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warning/0002/quant-ph0002009.html | ar5iv | text | # Information - The Fundamental Notion of Quantum Theory
## 1 Introduction
Recently, the experiments performed by the groups of Zeilinger and Weinfurter during the last years indicate that in understanding quantum theory the concept of information may play an essential role. Zeilinger proposed in \[BZ99\] an information measure suitable for use in quantum physics.
Here I try to introduce an information measure taking into account two main points of the discussion about interpretation of quantum theory:
1. The contextuality of any quantum measurement.
2. The transition from quantum regime to classical regime taking place during a measurement.
Because of the significance of the first point I will concentrate on it. Already Niels Bohr stressed in his interpretation the importance of taking into account the interdependance of measuring apparatus and quantum object. He always stated that in every measurement also the apparatus has to be described exactly in order to account for the observable properties of the quantum object in question. Only the observers’ questions decide which of the quantum objects’ property gets a precise well-defined value. Other properties related to it by an uncertainty relation, however, are still indetermined
## 2 Some basics
In this section I shortly describe some features of the density matrix formalism as far as necessary for my purpose. I restrict to the case of finite dimensional Hilbert space.
Two different types of quantum objects can be distinguished, which I classify according to the goal, the description of measurement in terms of information. The key notion is the description of quantum objects by a density matrix.
### 2.1 Quantum Objects in Pure State
I denote a single quantum object, unknown to the environment by the term *type 1-system*; such an object has never been measured. It likewise could be described by a $`\psi `$-function, written as $`\psi =_{i=1}^na_i\phi _i`$ , { $`\phi _i`$} an orthonormal basis of the underlying Hilbert space. Then its density matrix $`\widehat{\rho }`$ has the form:
$$\rho =\left(\begin{array}{ccccc}\left|a_1\right|^2& a_1a_2^{}& a_1a_3^{}& \mathrm{}& \mathrm{}\\ a_2a_1^{}& \left|a_2\right|^2& a_2a_3^{}& \mathrm{}& \mathrm{}\\ a_3a_1^{}& a_3a_2^{}& \left|a_3\right|^2& \mathrm{}& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\end{array}\right)\text{with}tr\rho =\underset{i=1}{\overset{N}{}}\left|a_i\right|^2=1$$
(1)
Such a hermitian matrix can be brought to diagonal shape by an unitary transformation to e.g.:
$$\widehat{\rho ^T}=\left(\begin{array}{ccc}1& 0& ..\\ 0& 0& ..\\ ..& ..& ..\end{array}\right)$$
(2)
This transformation could in principle be written as $`U=e^{ıHt}`$ with a suitable Hamiltonoperator $`H`$ and a suitable time $`t`$. But since the original density matrix is unknown, also the transformation can not be known in advance. There only could be an educated guess in order to achieve this form of the density matrix . This technique is used in the development of quantum computing (Grovers algorithm)(see e.g.\[CGK98\]) .
#### 2.1.1 Contextuality of quantum objects
The above matrix $`\rho `$ represents a pure state, i.e. $`\rho ^2=\rho `$. This property is invariant under unitary transformations. But the coefficients in the matrix representation depend on the basis chosen. We imagine that the state $`\psi `$ resp. $`\rho `$ is represented with respect to a specified measuring apparatus which means selection of a measurement basis $`\left\{\phi _i\right\}.`$ Hence the entries of the matrix reflect the relation between the state and the chosen measurement, or in other words, the quantum object in its context. If the density matrix looks like $`\widehat{\rho ^T}`$ (see (2)), then the state is an eigenstate with respect to the measurement basis. This can be interpreted as the quantum object being in a definite state relative to the corresponding measurement, i.e. the corresponding eigenvalue is attained with probability 1. If this is not the case the density matrix will be of shape (1) with at least two of the $`a_i0`$.
##### Example:
The simplest possible example is a single photon with density matrix
$$\rho _z=\left(\begin{array}{cc}a_1a_1^{}& a_1a_2^{}\\ a_1^{}a_2& a_2a_2^{}\end{array}\right)$$
relative to the basis formed by the eigenstates of the $`\sigma _z`$-Operator, say. It would look different with respect to the eigenstates of the $`\sigma _x`$-Operator, namely:
$$\rho _x=\frac{1}{2}\left(\begin{array}{cc}(a_1+a_2)(a_1+a_2)^{}& (a_1+a_2)(a_1^{}a_2^{})\\ (a_1+a_2)^{}(a_1a_2)& (a_1a_2)(a_1a_2)^{}\end{array}\right)$$
### 2.2 Ensemble of Quantum Objects
With the term *type 2-system* I denote an ensemble of quantum objects. The density matrix of such an ensemble is described by diagonal entries giving the probabilities of the corresponding measurement results (\[Hun96\]) and hence might be written as
$$\stackrel{~}{\rho }=\left(\begin{array}{ccccc}\left|a_1\right|^2& 0& 0& 0& \mathrm{}\\ 0& \left|a_2\right|^2& 0& 0& \mathrm{}\\ 0& 0& \left|a_3\right|^2& 0& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\end{array}\right)\text{with}tr\rho =\underset{i=1}{\overset{N}{}}\left|a_i\right|^2=1$$
(3)
In this case the behaviour of the quantum objects allows for an ignorance interpretation, where the properties of a single object are fixed but unknown. Therefore it could be modelled by throwing dice.
These two cases can be distinguished by a proper notion of quantum information.
## 3 Notion of Quantum Information
### 3.1 Quantum Information
The quantum information, i.e. the information present in a quantum object, will be defined as
$$I_Q:=Ctr\rho ^2C\text{a\hspace{0.17em} suitably\hspace{0.17em} chosen\hspace{0.17em} constant}$$
(4)
This defintion is motivated by the fact, that the density matrix incorporates all the properties of the object in question. The expectation value of any observable $`O`$ is given by: $`O:=tr(\rho O)`$. Hence the expression $`tr`$$`\rho ^2`$ can be considered as the expectation value of the information inherent in the quantum object. The normalization constant $`C`$ should be chosen as $`\mathrm{log}N`$, where $`\mathrm{log}`$ denotes the logarithm of basis 2 and $`N`$ is the dimension of the underlying Hilbert space. In this formulation the value $`C`$ corresponds to the whole available information, counted in bit or put in other words, the minimal number of questions necessary for determining the state of the quantum object. The value 0 - that cannot be attained - would correspond to the (impossible) case that there is no (quantum) information at all.
#### 3.1.1 Type 1 Systems
From the definition it is clear that $`I_Q=C=\mathrm{log}N`$ for every quantum object in a pure state (since $`\rho ^2=\rho `$). It is interesting to give a meaning to this result: A quantum object is, in strength, completely isolated from its environment. Hence this result can be interpreted in a way that the quantum object has the whole information - including its quantum mechanical surplus-knowledge (see below, 3.3) - incorporated on its own. And this internal information of the quantum object is independent of all other things that might happen in the world. Even changes internal to the quantum object do not have any influence on the amount of information as long as they correspond to unitary transformations. The quantum information is - by the very definition - always equal to $`C`$, i.e. always complete and always the whole information thinkable of.
#### 3.1.2 Type 2-systems
In the case of type 2-systems, however, $`I_Q<C`$ in general. Here the definition of quantum information gives - with respect to a suitable basis
$$I_Q=Ctr\stackrel{~}{\rho }^2=C\underset{i}{}\left|a_i\right|^4>0$$
where $`\left|a_i\right|^2`$are the diagonal elements of the density matrix $`\stackrel{~}{\rho }`$ (see (3)). This measure of information attains a minimum $`\frac{\mathrm{log}N}{N}`$ if all states are equally probable, i.e. $`\left|a_i\right|^2=\frac{1}{N}`$ for all $`i`$, and a maximum, namely $`\mathrm{log}N`$, if one state is attained with probability 1.
The difference between the two types of quantum systems hence is clearly visible on the basis of the notion of quantum information. This instrument can be more refined.
### 3.2 Interaction of quantum objects
Let two quantum objects - the object S and the object M with density matrices $`\rho _S`$; $`\rho _M`$ respectively - interact with each other. In this interaction case the definition
$$I_Q^I:=(C_M+C_S)tr(\rho _S\rho _M)^2$$
gives the information, object M carries about object S, say. $`I_Q^I`$ , by definition, does not change during the interaction as long as it is described by a Schrödinger equation.
We can distinguish three cases:
1. Both objects, S and M, are of type 1. Then the compound system is again of type 1, and an isolated quantum object with quantum information
$$I_Q^I=(C_M+C_S)tr(\rho _S\rho _M)^2=C_M+C_S$$
This means, in the context of two quantum objects, that both carry the full information of *each* *other* because of the entanglements arising between them because of the interaction.
2. Both objects, S and M, are of type 2. Then the compound system is again of type 2 with quantum information
$$I_Q^I=(C_M+C_S)tr(\rho _S\rho _M)^2=(C_M+C_S)\underset{i}{}\left|a_i^S\right|^4\underset{j}{}\left|a_j^M\right|^4=\frac{(C_M+C_S)}{C_SC_M}I_Q^SI_Q^M$$
where $`I_Q^S;I_Q^M`$ are the quantum information of object S and object M, respectively.
3. The third case is the most interesting case because it can be used for a characterization of measurement: The object S is of type 1 and the object M of type 2. We have:
$$I_Q^I=(C_M+C_S)tr(\rho _S\rho _M)^2=(C_M+C_S)\frac{I_Q^S}{C_S}\frac{I_Q^M}{C_M}=I_Q^M(1+\frac{C_S}{C_M})<C_M+C_S$$
The strict inequality indicates that the measuring object M takes information from object S, but in general not the whole information $`I_Q^S=C_S`$. (Of course M still holds its “own” information $`I_Q^M`$.)
### 3.3 The quantum mechanical surplus-knowledge
The information $`I_Q`$ only “sees” the quantum object, not any relation to a measurement. Its constant value $`C`$ for a pure state reflects the fact that a quantum object always carries the whole information about its state in it. In a measurement, however, only parts of this information come into “reality”. The other parts are called the “quantum mechanical surplus-knowledge” by Weizsäcker, \[vW94\], and Görnitz \[Gör99\] stresses the importance of the relations between different parts of a quantum object. Hence the off-diagonal elements of the density matrix seem to be an appropriate measure for this “surplus knowledge”. As alluded to before, (see section 2.1), the off-diagonal elements depend on the kind of contact with environment (measurement) or, in other words, on the relation between the state of the quantum object and the (planned) measurement. This observation also reflects the considerations of Bohr who always stressed that the appearance of a quantum object depends on the kind of measurement. The surplus-knowledge hence is deeply connected to the basis chosen, i.e. to the “planned measurement”.
How to define the surplus-knowledge? Let $`\rho `$ be the density matrix of a quantum object S of shape (1) and $`\stackrel{~}{\rho }`$ the corresponding diagonal matrix $`\stackrel{~}{\rho }=diag_i(\left|a_i\right|^2)`$.
The relation between $`\rho `$ and $`\stackrel{~}{\rho }`$ can be interpreted in a twofold way:
1. Given a quantum object S with density matrix $`\rho `$ we get $`\stackrel{~}{\rho }`$ by a complete measurement, in the end, equivalent to the density matrix of a type 2-system or an ensemble.
2. Or vice versa, given $`\stackrel{~}{\rho }`$, - the density matrix belonging to an ensemble - we reconstruct the state $`\rho `$ (1) of the quantum object S from the diagonal elements of $`\stackrel{~}{\rho }`$.
Let us now define the off-diagonal information, the “surplus-knowledge” contained in the density matrix $`\rho `$ of quantum object S, as
$$K_Q^S:=Ctr(\rho \stackrel{~}{\rho })^2$$
(5)
$`K_Q^S`$ can be expressed, as desired, in terms of the off-diagonal elements of the density matrix $`\rho `$:
$$K_Q^S=Ctr(\rho \stackrel{~}{\rho })^2=C\underset{ij}{}\left|a_ia_j^{}\right|^2$$
resp.
$$K_Q^S=C\underset{i}{}\left|a_i\right|^2(1\left|a_i\right|^2)=CC\underset{i}{}\left|a_i\right|^4=C\stackrel{~}{I_Q}\text{where}\stackrel{~}{I_Q}=Ctr\stackrel{~}{\rho }^2$$
(6)
This expression admits two interpretations:
1. $`K_Q^S`$ may be interpreted as the difference between the information obtained from the ontological and from the epistemical interpretation of a quantum object, because $`\stackrel{~}{I_Q}=Ctr\stackrel{~}{\rho }^2`$ reflects the epistemological knowledge contained in the quantum object in question. Since $`\stackrel{~}{\rho }`$ is diagonal we furthermore have $`Ctr(\rho \stackrel{~}{\rho })^2=Ctr\rho ^2Ctr\stackrel{~}{\rho }^2`$.
Hence, the whole information $`I_Q`$ of a quantum object can be divided into a classical part - contained in the diagonal elements - and a quantum part \- contained in the off-diagonal elements, i.e. $`I_Q=Ctr\rho ^2=C=K_Q^S+\stackrel{~}{I_Q}`$. I again want to stress that the quantum part of the information depends on the relation of state and measurement, i.e. the measured observable. There *always* is a measurement relative to which a quantum object is in a *pure* state (2.1). But simultaneously it is undetermined with respect to non-commuting observables (s.a. the example in 2.1.1). In the first case there is no surplus knowledge, $`K_Q^S=0`$, (relative to the fixed measurement, which with probability 1 shows a fixed value for the measured observable), but in the second case $`K_Q^S0`$. Hence the occurence of a non-vanishing surplus-knowledge is deeply connected to the uncertainty relations.
2. In a second interpretation the surplus-knowledge $`K_Q^S`$ can be regarded as the possible information gain during a measurement or the information exchange between the quantum object and its environment (resp. measuring apparatus): If $`\rho `$ describes an (unknown) object before measurement and $`\stackrel{~}{\rho }`$ the (partly) known object after a measurement then $`Ctr\rho ^2`$ is the information contained in the unmeasured quantum object (normally equal to $`C`$) and $`\stackrel{~}{I_Q}=Ctr\stackrel{~}{\rho }^2`$ is the information still contained in the quantum object after measurement. This would correspond to building a partial trace in the standard density matrx formalism.
## 4 Working with the notion of information
The sense and function of these notion can best be explored at work.
### 4.1 Examples
As the simplest possible example we treat the case of one resp. two photons.
#### 4.1.1 Case of single photon
A single photon can be written as $`\psi _1=a_1|0+a_2|1`$ with $`\left|a_1\right|^2+\left|a_2\right|^2=1`$. This corresponds \- relative to the standard basis $`|0=\left(\begin{array}{c}1\\ 0\end{array}\right);|1=\left(\begin{array}{c}0\\ 1\end{array}\right)`$ \- to a density matrix
$$\rho _1=\left(\begin{array}{cc}a_1a_1^{}& a_1a_2^{}\\ a_1^{}a_2& a_2a_2^{}\end{array}\right)$$
The quantum information is $`I_Q=1`$ and the quantum mechanical surplus-knowledge then is $`K_Q^S=2\left|a_1a_2\right|^2\frac{1}{2}`$. It is determined in relation to a spin measurement along the directions $`|0;|1`$. Relative to the representation of $`\psi `$ in the basis $`a_1|0+a_2|1,`$ $`a_1|0a_2|1`$ the density matrix would look like: $`\rho _1=\left(\begin{array}{cc}1& 0\\ 0& 0\end{array}\right)`$ with no surplus knowledge, because the photon then is in a eigenstate relative to the corresponding measurement. The relation between the surplus-knowledge and the whole quantum information marks the amount of information extractable from the quantum object - in a fixed context.
#### 4.1.2 Product state of two photons
Let us assume that another photon $`\psi _2=b_1|0+b_2|1`$ is brought into contact with the first photon. This can result in a type 1-system which means that both photons get entangled and give rise to the most general density matrix
$$\rho =\rho _1\rho _2=\left(\begin{array}{cccc}\left|a_1b_1\right|^2& \left|a_1\right|^2b_1b_2^{}& a_1a_2^{}\left|b_1\right|^2& a_1a_2^{}b_1b_2^2\\ \left|a_1\right|^2b_1^{}b_2& \left|a_1b_2\right|^2& a_1a_2^{}b_1^{}b_2& a_1a_2^{}\left|b_2\right|^2\\ a_1^{}a_2\left|b_1\right|^2& a_1^{}a_2b_1b_2^{}& \left|a_2b_1\right|^2& \left|a_2\right|^2b_1b_2^{}\\ a_1^{}a_2b_1^{}b_2& a_1^{}a_2\left|b_2\right|^2& \left|a_2\right|^2b_1^{}b_2& \left|a_2b_2\right|^2\end{array}\right)$$
which already is properly normalized with quantum information $`I_Q=2`$ and surplus knowledge
$`K_Q^S=4(\left|a_1\right|^4\left|b_1b_2\right|^2+\left|a_1a_2\right|^2\left|b_1\right|^4+2\left|a_1a_2b_1b_2\right|^2+\left|a_1a_2\right|^2\left|b_2\right|^4+\left|a_2\right|^4\left|b_1b_2\right|^2)`$
$`=4((\left|a_1b_1\right|^2+\left|a_2b_2\right|^2)(\left|a_1b_2\right|^2+\left|a_2b_1\right|^2)+2\left|a_1a_2b_1b_2\right|^2)\frac{3}{2}`$.
If one of the photons would be in a eigenstate (e.g. $`b_1=1;b_2=0`$) this surplus knowledge would reduce to $`4\left|a_1a_2\right|^21`$.
In general a system of $`n`$ 2-state quantum objects with equal probabilities $`\frac{1}{2}`$ for all outcomes of a measurement possesses the (maximal possible) surplus-knowledge
$$K_{Q,max}^S=n(1\frac{1}{2^n})$$
#### 4.1.3 EPR-pairs of photons
EPR-pairs are of special interest. Their density matrix can not be written in terms of the product of the density matrices of the single potons (this is the way they are constructed). Let us assume the singlett state $`\psi =\frac{1}{\sqrt{2}}(|0,1|1,0)`$. Herewith
$$\rho =\frac{1}{2}\left(\begin{array}{cccc}0& 0& 0& 0\\ 0& 1& 1& 0\\ 0& 1& 1& 0\\ 0& 0& 0& 0\end{array}\right)$$
with respect to the basis built by eigenstates of the operator $`\sigma _z`$. This gives $`I_Q^{\text{EPR}}=2`$ and the *maximal possible* surplus-knowledge $`K_Q^{\text{EPR}}=1`$. One *single* measurement can make the whole system “classical” i.e. well determined with respect to this fixed measurement, the $`\sigma _z`$\- observable (see below section Criterion:).
The same is valid for the so-called GHZ-states which have a similar density matrix in the non-zero parts of their density matrix. Correspondingly $`I_Q^{\text{GHZ}}=3`$ and $`K_Q^{\text{GHZ}}=\frac{3}{2}`$. As seen below, (sections Criterion:, 4.2.1) also this system becomes “classical” with respect to a fixed measurement observable in a single measurement.
#### 4.1.4 Ensemble of identical photons
The density matrix of an ensemble of $`n`$ identical photons is built as the sum of the single density matrices. For further analysis I separately introduce arbitrary phases $`\phi _i`$ such that
$$\rho _{i,single}=\left(\begin{array}{cc}a_1^2& a_1a_2e^{ı\phi _i}\\ a_1a_2e^{ı\phi _i}& a_2^2\end{array}\right)$$
Hence
$$\rho _{system}=\frac{1}{n}\underset{i=1}{\overset{n}{}}\rho _{i,single}=\frac{1}{n}\left(\begin{array}{cc}na_1^2& a_1a_2_{i=1}^ne^{ı\phi _i}\\ a_1a_2_{i=1}^ne^{ı\phi _i}& na_2^2\end{array}\right)$$
The constant $`C`$ is exactly $`1`$ in this case. Hence there is the surplus knowledge
$`K_Q^S(\rho _{system})={\displaystyle \frac{1}{n^2}}\left[\left|{\displaystyle \underset{i}{}}a_1a_2e^{ı\phi _i}\right|^2+\left|{\displaystyle \underset{i}{}}a_1a_2e^{ı\phi _i}\right|^2\right]`$
$`={\displaystyle \frac{\left|a_1a_2\right|^2}{n^2}}\left[\left|{\displaystyle \underset{i}{}}e^{ı\phi _i}\right|^2+\left|{\displaystyle \underset{i}{}}e^{ı\phi _i}\right|^2\right]`$
and the quantum information
$$I_Q=(a_1^4+a_2^4)+\frac{a_1^2a_2^2}{n^2}\underset{i,j=1}{\overset{n}{}}e^{ı(\phi _i\phi _j)}=(a_1^4+a_2^4)+\frac{2a_1^2a_2^2}{n^2}\underset{i,j=1}{\overset{n}{}}\mathrm{cos}(\phi _i\phi _j)$$
Now two extreme cases can be distinguished:
1. We assume $`n`$ is very large and the phases $`\phi _i`$ are distributed randomly with equal weight. Then the term containing the diagonal elements outweighs the other term depending on the phases. Hence $`I_Q(a_1^4+a_2^4)<1`$ and $`K_Q^S0`$. This indicates a transition from a quantum system to a “classical” system where the non-knowing of measurement results can be interpreted epistemically.
2. We assume all the phases $`\phi _i`$ are equal to a single phase $`\phi `$ . Then the expressions for the informations simplify to $`K_Q^S=2\left|a_1a_2\right|^2`$, the surplus knowledge contained in a single photon, and to $`I_Q=(a_1^2+a_2^2)^2=1`$. Taken together this indicates that the ensemble constitutes a quantum object with only two possible states, corresponding to an ensemble of coherent photons, behaving like *one* single photon. In this place it is quite interesting to note that hence an ensemble of identical photons as required e.g. in the ensemble interpretation, cannot be distinguished from a single photon. Both carry the same information content and the same surplus-knowledge.
For convenience I give the formulas for a system of two identical photons which already display all the described behaviour: the density matrix is
$$\rho _{system}=\frac{1}{2}\left(\begin{array}{cc}2a_1^2& a_1a_2(1+e^{ı\phi })\\ a_1a_2(1+e^{ı\phi })& 2a_2^2\end{array}\right)$$
with $`K_Q^S=\frac{\left|a_1a_2\right|}{2}^2(1+\mathrm{cos}\phi )`$ and $`I_Q=(a_1^4+a_2^4)+\frac{\left|a_1a_2\right|^2}{2}(\mathrm{cos}\phi +1)`$, where $`\phi `$ is the relative phase of the two photons. If $`\phi =0`$, then the density matrix of the whole system $`\rho _{\text{system}}`$ is just that of the single photon’s density matrix $`\rho _{\text{single}}`$. This again shows that two identical photons together are described by *one* wavefunction, exhibiting single particle behaviour. A similar phenomenon occurs in the Bose-Einstein-condensates.
### 4.2 Quantum Information and Measurement
With aid of these notions of quantum information and surplus knowledge we now approach the measurement problem. We do not go into any detailed discussion of the measurement problem; this can be read elsewhere \[BLM91\].
#### 4.2.1 Assumptions for measurement
By definition we cannot know anything about pure quantum objects, i.e. objects of type 1. Hence we make
Only quantum objects of type 2 can be used as a measuring apparatus (see also \[BLM91\]).
By the very definition we can have knowledge only about quantum objects of type 2, because they allow for an epistemical interpretation of quantum objects; i.e. there are fixed values for the properties to be measured, the observer only does not know which value is realized. Objects of type 1 in contrast do not have fixed values for its properties at all; the properties of those objects come into existence only with a measurement.
Hence the measurement problem is most deeply related to interaction case 3 from section 3.2.
Now we assume quantum objects with density matrices as in 1 and in 3 and using the same notation we state the reduction postulate as:
If after the measurement any $`a_j^M=0`$ (that has been different from 0 before measurement), then there are at least one index $`i_1,\mathrm{}.,i_r`$ such that $`a_{i_1}^S=\mathrm{}=a_{i_r}^S=0`$.
This assumption goes just the other way round than most other assumptions on measurement devices. In my opinion this formulation gives the possibility of dealing with the phenomenon of so-called quantum erasers.
A measurement is only fruitful and hands over new information from the quantum object to the classical regime if more than one of the $`a_i^M`$or $`a_i^S`$ are different from zero. A measuring apparatus $`M`$ should give a statement which allows to draw conclusions on the quantum object $`S`$. I.e. if a possible result of $`M`$ is excluded with probability one (that is one $`a_j^M=0`$), then there should be properties of $`S`$ that also can be excluded with probability one. This seems to me to be a reasonable assumption because otherwise any measurement would be completely useless or put differently: the result $`a_j^M=0`$ of $`M`$ would give no information on $`S`$, i.e. it would be no true measurement.
As we have seen, the information that can be extracted from a quantum object depends on the design of the measurement or - more generally - on the environment it is brought into contact with. Furthermore - in order to extract and interpret the information - we have to know something about the measuring device.
#### 4.2.2 Criterion for completion of measurement
Given a state and a fixed measurement (observable) the quantum mechanical surplus-knowledge $`K_Q^S`$ measures the quantum object’s degree of being “quantum” with respect to this measurement. From the definition of $`K_Q^S`$ we can define a quantum object as “behaving classical” if $`K_Q^S`$ is sufficiently small. Classically there cannot be an amount of information less than 1 bit; so we set:
The object S can be regarded as a classical object with respect to a fixed measurement if the corresponding surplus knowledge $`K_Q^S<1.`$ Then we regard the measurement as completed (at least as far as possible).
This condition is always fulfilled in case of a 2-dimensional Hilbert space, i.e. e.g. for the spin states of a single photon. But already for 3-dimensional Hilbert space, there has to be some measurement in order to reduce the quantum mechanical surplus knowledge before a quantum object can be regarded as “classical”, i.e. as having fixed values for one property.
### 4.3 Example for Measurement: Photon Interference
Now let us treat the most important case, the case of a photon being measured by a object showing two possibilities. This setting is given for example by the double slit experiment, a Michelson interferometer or similar devices.
Of special interest in this context is the treatment of the phenomenon of quantum eraser and which-way-information (\[HKWZ96\]).
As alluded to before an interference experiment would correspond to case 3: the interaction of the photon with a measuring device having two well defined possibilities (ways).
##### Before the measurement:
The density matrix of an interferometer (before the measurement) is given by:
$$\rho _M=\left(\begin{array}{cc}\alpha ^2& 0\\ 0& \beta ^2\end{array}\right)\text{with}\alpha ^2+\beta ^2=1$$
and the density matrix of the photon is as described above (see (1)). Then the common density matrix of photon and measuring device is given by
$$\rho _S\rho _M=\left(\begin{array}{cccc}\alpha ^2a_1a_1^{}& \alpha ^2a_1a_2^{}& 0& 0\\ \alpha ^2a_1^{}a_2& \alpha ^2a_2a_2^{}& 0& 0\\ 0& 0& \beta ^2a_1a_1^{}& \beta ^2a_1a_2^{}\\ 0& 0& \beta ^2a_1^{}a_2& \beta ^2a_2a_2^{}\end{array}\right)$$
The definition of quantum information gives:
$$I_Q^I:=2tr(\rho _S\rho _M)^2=2I_Q^M=2(\alpha ^4+\beta ^4)$$
We distinguish two cases:
1. $`\alpha ,\beta 0`$. In this case $`I_Q^I<2`$, where the strict inequality reflects the fact that the compound system is not a whole quantum object but includes a “semiclassical” device. The smaller $`I_Q^I`$ the more the measurement device makes the quantum object “classical”. On the other hand the quantum object can not be interpreted completely as classical as is seen from the surplus-knowledge $`K_Q^I:=2_{ij}\text{off-diagonal}^2=4(\alpha ^4+\beta ^4)\left|a_1a_2\right|^2=I_Q^IK_Q^S`$, indicating the quantum character of the compound system. The minimum of both, $`I_Q^I`$ and $`K_Q^I`$, therefore is attained if both “ways” can be discerned clearly, i.e. $`\alpha ^2=\beta ^2=\frac{1}{2}`$. Then $`K_Q^I=2\left|a_1a_2\right|^2=K_Q^S`$, the surplus knowledge of the single photon which only vanishes if $`a_1=0`$ or $`a_2=0`$ meaning that the quantum object would be in a eigenstate with respect to the chosen measurement.
2. $`\alpha =1;\beta =0`$ or vice versa. In this case the measurement can give no information on the quantum object to the outside environment and the quantum information remains undisturbed. Here the surplus knowledge $`K_Q^I`$ becomes maximal - $`K_Q^I=4\left|a_1a_2\right|^2`$ \- corresponding to the fact, that - no matter what is done during the “measurement” - we cannot distinguish different states of the quantum object in question. That means there is no true measurement, because no information is extracted from the quantum object.
These information values characterize the compound system only before a measurement: the quantum object interacts with the measurement device without being read from the outside.
##### After the measurement:
After the measurement of a single photon, we should know “which way it has taken”, simulating the situation as if (for this single photon) $`\alpha =1`$resp. $`\beta =0`$ and from this information we would like to draw conclusions on the quantum object. A measurement on an ensemble requires that for every single photon it has to be decided whether “the photon has taken way 1” or whether “the photon has taken way 2”, giving in the end the respective probabilities $`\alpha ^2`$ resp. $`\beta ^2`$. Hence from (the factual) $`\beta `$ equal to 0 for a single photon we should be able to conclude e.g. (the fact) $`a_1=0`$, according to assumption 2 (reduction postulate, 4.2.1). Then $`I_Q^{I,\text{after}}=2`$, that is: M carries the whole information of the interaction and $`K_Q^{S,\text{after}}=4\left|a_1a_2\right|^2=0`$. Then the measurement is completed and the photon may be regarded as a classical object with definite properties *\- but!: with respect to the performed measurement only!*
But since $`I_Q^{I,\text{after}}`$ has attained its maximal possible value, there might still be non-vanishing surplus knowledge with respect to other measurements (observables).
If, as for instance in a quantum eraser, assumption 2 (see 4.2.1) is hurt, i.e. $`\beta =0`$ *and* $`a_10`$ after a measurement, $`K_Q^{S,\text{after}}`$ is different from 0, i.e. the photon is *not* in a eigenstate with respect to the basis of the measurement or, in other words, (part of) the information is left to the quantum object.
## 5 Conclusion
In this article there were introduced two notions of quantum information reflecting the differences between quantum and classical objects. These give the fundamental notion of information a quantitative expression and show clearly the contextuality of quantum objects.
I have to thank Prof. Thomas Görnitz for innumerable helpful discussions and many valuable hints. |
warning/0002/hep-th0002174.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Brane world scenarios are based on the hypothesis that the three space dimensional world that we appear to be living in is a brane that is embedded in a higher dimensional world<sup>1</sup><sup>1</sup>1This is an old idea that has been revived recently in a string inspired context in .. Most of the work on this has been of a phenomenological nature and not many attempts have been made to justify the postulates within a well defined framework for (higher dimensional) quantum gravity such as string theory. Nevertheless this activity is “string inspired”, in that an obvious candidate for such a world is a collection of D-branes on which (at least in principle) the standard model can live. For most of this paper we will not worry about a string realization though towards the end we will suggest some possibilities.
The main issue that we are concerned with here, is that of obtaining flat 3+1 dimensional solutions to the equations of the effective higher dimensional theory in a natural way (i.e. without fine tuning). We will show that there are situations where flat brane solutions can be obtained by choice of integration constant<sup>2</sup><sup>2</sup>2After this work was substantially completed several papers appeared which obtain a one flat brane solution by choice of integration constant . We will comment on these works in the conclusion.. In this respect this mechanism is a realization whithin the brane world context of an old idea going back to ,,,,<sup>3</sup><sup>3</sup>3There are also unimodular gravity scenarios that appear to have been first discussed in and are also discussed in . However these do not fit in naturally in a string picture.. To set the stage for the brane world calculations we will first review this argument.
Consider an effective theory describing our four dimensional world at low energies of the form
$$S=\frac{1}{2\kappa ^2}\left(\sqrt{G}R\frac{1}{2}F_4^{}F_4\right)+S_m(G,\psi ),$$
(1)
where $`S_m`$ is the matter action and $`F_4`$ is a four form field strength satisfying the Bianchi identity $`dF_4=0`$. The equations of motion from this action are,
$`R_{\mu \nu }{\displaystyle \frac{1}{2}}G_{\mu \nu }R{\displaystyle \frac{1}{2.4!}}(4F_\mu \mathrm{}F_\nu ^{\mathrm{}}{\displaystyle \frac{1}{2}}G_{\mu \nu }F^2)`$ $`=`$ $`2\kappa ^2T_{\mu \nu }`$
$`d^{}F`$ $`=`$ $`0.`$ (2)
In the above $`T_{\mu \nu }`$ is the matter stress tensor and we have ignored the matter equations of motion which will not play any role in this paper. Now the four form equation of motion and Bianchi identity have the solution,
$$F_4=\mu ^{}1$$
(3)
where $`\mu `$ is a constant and the second factor is the volume form. When this is substituted into the first equation one gets
$$R_{\mu \nu }=G_{\mu \nu }(2\kappa ^2V_0+\frac{\mu ^2}{4})$$
(4)
Here $`V_0`$ is the effective cosmological constant generated in the matter sector. Clearly if this is negative then the integration constant $`\mu `$ can always be chosen so as to get flat space. The question is what is the significance of this result. We should empahsize here that from the perspective of a four dimensional theory the integration constant $`\mu `$ can take any real value. However when the theory is embedded in a higher dimensional theory such as string theory which admits branes which are sources of the four form field, $`\mu `$ is quantized and this mechanism would not work. We will encounter such a situation later on when we discuss what happens in string theoretic brane world scenarios.
First note that if one wants to argue that a flat space solution can be obtained, even in the presence of quantum corrections to the matter action (ignoring gravity sector fluctuations), then one should replace the classical matter action $`S_m`$ by the quantum effective action $`\mathrm{\Gamma }_m`$. $`V_0`$ is now dependent on the RG scale and so the integration constant needs to be renormalization scale dependent in order to get flat space at every scale. Of course such a constant can be chosen at will but to solve the cosmological constant problem the question of why out of the real line of values of this integration constant, one particular value (or one value at each scale) gets chosen should also be answered. At present there is no clear answer to this and we will discuss this question further in section IV.
Nevertheless one may take the point of view that replacing a fine tuning problem with a choice of integration constant is progress, since one is not adjusting a parameter in the Lagrangian. In fact in string theory there are no parameters to adjust and one might well need a mechanism like this to get flat space solutions after supersymmetry breaking. So it might still be worth investigating whether such mechanisms are available there.
In the next section we will motivate a brane world scenario from a bottom up approach as opposed to a top down string approach by asking whether the RG scale in four dimensions can be thought of as a fifth dimension. In section three we will discuss explicit embeddings of branes in five dimensions and discuss how the flat one and two brane solutions emerge. In fact in the string theory case we will argue that the natural scenario is a two-brane one. From the five dimensional point of view this requires a fine tuning of a parameter in the bulk potential, but we will argue that there are compactifications of string theory in which this parameter is (from the ten dimensional point of view) an integration constant. In the concluding section we discuss the problem of justifying the choice of integration constants that leads to flat branes.
## 2 Renormalization Group Flow in External Gravity.
Let us consider the quantum theory corresponding to the classical action $`S_m`$. The fields $`\psi `$ could stand for the full set of standard model fields and we will also include a dilaton $`\varphi `$ in order to make the connection later on to string theory. We are going to do semi-classical dilaton-gravity. In other words the dilaton gravity sector is treated classically while the standard model fields are treated quantum mechanically. The quantum theory is defined by the functional integral,
$$e^{iW[G,\varphi ]}=[d\psi ]e^{iS_m[G,\varphi ,\psi ]}.$$
(5)
In order to define the quantum theory in a general gravitational background, a proper time cutoff propagator
$$K_ϵ^1=_{ϵ_0}^ϵe^{Ks}𝑑s$$
(6)
is introduced with $`K`$ being the kinetic operator. Here $`ϵ_0`$ may be regarded as the ultra-violet cutoff (taken for instance to be the string scale) and $`ϵ`$ may either regarded as a renormalization scale or the scale defining a Wilsonian effective action. Using also the technique of Riemann normal coordinate expansions, one can derive in principle the quantum effective action in a systematic way preserving general covariance. The quantum action can therefore be written in a derivative expansion as
$$W[G,\varphi ]=d^4x\sqrt{G}(\mathrm{\Phi }(\varphi ,ϵ)RZ(\varphi ,ϵ)(\varphi )^2+V(\varphi )+\mathrm{}$$
(7)
where the elipses represent higher derivative terms. We have indicated the explicit dependence on the RG scale. There would also of course be implicit dependence since the external fields $`G,\varphi `$,like the couplings of the theory will aquire $`ϵ`$ dependence. Also we have set all expectation values of standard model fields to their values solving the equations of motion (at this point the functional $`W`$ is in fact equal to the 1PI effective action $`\mathrm{\Gamma }`$) and have been suppressed. The RG equation reads,
$$\frac{dW}{dϵ}=\frac{W}{ϵ}+\beta _\lambda \frac{W}{\lambda }+\beta _{\mu \nu }.\frac{\delta W}{\delta G_{\mu \nu }}+\beta _\varphi .\frac{\delta W}{\delta \varphi }=0$$
(8)
where the $`\lambda `$ are the couplings in the theory with associated beta function $`\beta _\lambda `$ and the other betas are the analogous beta functions for the metric and phi field (which are to be treated as generalized couplings). When the classical action for gravity and the $`F_4`$ field are added to the above quantum action we again get an action of the form of (1) (plus higher derivative terms) but with couplings which depend on $`\varphi `$ and the RG scale $`ϵ`$. After a Weyl transformation this can be written as
$$S=\frac{1}{2\kappa ^2(ϵ)}\left(\sqrt{G}R\stackrel{~}{Z}(\varphi ,ϵ)(\varphi )^2\frac{1}{2}U(\varphi ,ϵ)F_4^{}F_42\kappa ^2V(\varphi ,ϵ)+\mathrm{}\right)$$
(9)
The previous argument still goes through with slight modifications. For instance now the four form equation is replaced by $`d^{}(U(\varphi )F_4)=0`$ which is solved by
$$F_4=\mu U^11$$
(10)
(which also satisfies the Bianchi identity). But the main result remains unchanged. The cosmological constant is an integration constant which can be chosen (in a RG scale dependent way) so as to get the effective cosmological constant to be zero. The argument is robust under renormalization of the standard model since it did not depend on particular functional forms of $`Z,u`$ or $`V`$. The problem of justifying the choice of integration constant however remains.
Let us now ask the question under what circumstances can the RG scale of the four dimensional theory be interpreted as a fifth dimension. In the argument was made that the five dimensional Hamilton-Jacobi equation can be interpreted as a four dimensional RG equation. Here we ask the opposite question; under what conditions can the latter be interpreted as a five dimensional gravity theory?
Consider the following expression constructed in terms of the quantum effective action $`W`$ defined in (7),
$`{\displaystyle \frac{1}{\sqrt{G}}}{\displaystyle \frac{1}{3}}\left(G^{\mu \nu }{\displaystyle \frac{\delta W}{\delta G^{\mu \nu }}}\right)^2`$ $``$ $`{\displaystyle \frac{\delta W}{\delta G^{\mu \nu }}}{\displaystyle \frac{\delta W}{\delta G_{\mu \nu }}}{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{\delta W}{\delta \varphi }}\right)^2`$ (11)
$`=`$ $`\sqrt{G}(\stackrel{~}{V}(\varphi ,ϵ)+{\displaystyle \frac{1}{\kappa ^2(\varphi ,ϵ)}}R+M(\varphi ,ϵ)(\varphi )^2+\mathrm{}`$
The right hand side is just a consequence of general covariance and the ellipses stand for higher derivative terms. The particular form of the expression on the left hand side is of course chosen to agree with the corresponding expresion in the Hamilton-Jacobi equation of five dimensions . Under what conditions can $`W`$ be interpreted as a classical five dimensional action? Clearly this is possible if the explicit dependence on $`ϵ`$ is absent.<sup>4</sup><sup>4</sup>4It should be noted that this explicit dependence includes the dependence on $`ϵ`$ through the renormalization of the flat space couplings as well. i.e. it corresponds to the first two terms of (8). It is possible that this is the case in $`𝒩=4`$ SU(N) Yang-Mills theory (at least in the large $`N`$ limit) and this would then be an explanation of the AdS/CFT conjecture .
Now the semi-classical theory of quantum fields is obtained after one adds a classical action and one then gets the action (9). Let us set the $`F_4`$ terms to zero for the moment and ask what happens to the cosmological constant. Let $`\varphi =\varphi _0(ϵ)`$ be a constant field satisfying $`\frac{V(\varphi ,ϵ)}{\varphi }=0`$. The gravity equation then gives $`R_{\mu \nu }=\frac{1}{2}V(\varphi _0(ϵ),ϵ)G_{\mu \nu }`$. Clearly if the explicit dependence of $`V`$ on $`ϵ`$ is absent then $`\varphi _0`$ is $`ϵ`$ independent and so is the Ricci curvature so that if one has tuned the minimum of $`V`$ to zero at some scale (for instance $`ϵ=ϵ_0`$) then one will get flat space at all scales. But the issue is precisely for what theories in four dimensions is the statement of independence from $`ϵ`$ valid. With sufficient supersymmetry this could be the case. But with $`𝒩=1`$ SUSY although the superpotential is not renormalized the Kahler potential is, so that the potential for $`\varphi `$ will in general depend explicitly on $`ϵ`$, though of course in this case one does not expect renormalization of the minimum of the potential. Thus in order to have a flat space solution at any RG scale one would in general need something like the mechanism discussed earlier.
Now it may be the case that, the absence of explicit dependence on $`ϵ`$ in $`W`$, while a sufficient condition for the five dimensional interpretation, is not be a necessary one. In other words there could be a cancellation of the epsilon dependence on the LHS of (11) amongst the different terms so that the RHS is $`ϵ`$ independent. In this case just the mere fact that a five dimensional interpretation (as in the AdS/CFT case ) exists, is no gurantee of RG invariance of the four dimensional cosmological constant<sup>5</sup><sup>5</sup>5Some discussion of the consequences of this are found in .. In other words the logic cannot be reversed. The absence of explicit dependence of $`W`$ on $`ϵ`$ (which implies in particular that the cosmological constant is not renormalized) is a sufficient condition for a five dimensional interpretation, but the latter does not imply that the former is the case.
## 3 Brane World Scenarios
In the previous section we discussed the assumptions that would lead us to interpret the RG scale as a fifth coordinate and thus four dimensional semi-classical gravity as a five dimensional gravity theory. Here we will explicitly treat the four dimensional theory as living on a brane in five dimensions. It is important to keep in mind the distinction between the two cases. In the first case the five dimensional theory (as for example in the AdS/CFT case) is simply a dual representation of the four dimensional quantum effective action. In the present case the underlying theory is five (or more dimensional) and the standard model is confined to a 3-brane living in it. This may perhaps be realized in string theory as for example a type IIB orientifold (compactified on some compact 5-manifold) with D3 branes and we shall discuss this further at the end of this section.
Using only general covariance, and keeping only two derivative terms, the most general five dimensional action of gravity coupled to a scalar field is,
$$S[G,\varphi ]=d^5x\sqrt{G}(RZ(\varphi ,ϵ)(\varphi )^2+V(\varphi )+\mathrm{})$$
(12)
If this originates from the string theory example mentioned above, the potential $`V`$ may come from the $`F_5`$ terms that occur there, just like the $`F_4`$ terms in equation (9), after using the solution to the equation of motion for the $`F_5`$ field<sup>6</sup><sup>6</sup>6Thus in the notation of (9) and the sentence below it, ( rewritten for five dimensions) $`V`$ in the above would be $`\frac{1}{2}\mu ^2U(\varphi )^1`$.
Let us take the coordinates to be $`x^M,M=0,1,\mathrm{..4}`$ with the fifth coordinate $`x^4=u`$. Now we insert 3-branes transverse to the direction $`u`$ at the points $`u=u_i`$. We choose the static gauge so that the embedding functions are $`x^\mu (\xi )=\xi ^\mu ,\mu =0,\mathrm{..3}`$ and ignore their fluctuations. The effective action(s) coming from integrating the “standard model” quantum fields (and hidden sector fields if there is more than one brane) will then take the form.
$$\underset{i}{}_{u=u_i}T_i(\varphi ,ϵ)\sqrt{G_{4(i)}}$$
(13)
There will also be derivative terms but since we are interested in solutions with flat metrics and constant fields in 4d, they are irrelevant to our discussion. The field equations for the system are then obtained by extremizing the sum of the two actions (12,13).
Now as in ,, we look for solutions that give flat space and constant $`\varphi `$ field on the brane. So we write
$`ds^2`$ $`=`$ $`e^{2\omega (u)}\eta _{\mu \nu }dx^\mu dx^\nu +du^2`$
$`\varphi `$ $`=`$ $`\varphi (u).`$ (14)
The equations of motion then take the following form, (writing $`\frac{d}{du}^{}`$)
$`6\omega ^2`$ $`=`$ $`{\displaystyle \frac{1}{2}}Z(\varphi )\varphi ^2V(\varphi )`$
$`3\omega ^{\prime \prime }+6\omega ^2`$ $`=`$ $`{\displaystyle \frac{1}{2}}[Z(\varphi )\varphi ^2+V(\varphi )]{\displaystyle \frac{1}{2}}{\displaystyle \underset{i}{}}2\kappa ^2T_i(\varphi ,ϵ)\delta (uu_i)`$
$`Z(\varphi ^{\prime \prime }+4\omega ^{}\varphi ^{})+{\displaystyle \frac{1}{2}}\varphi ^2{\displaystyle \frac{dZ}{d\varphi }}`$ $`=`$ $`{\displaystyle \frac{1}{2}}V^{}(\varphi )+\kappa ^2{\displaystyle \underset{i}{}}{\displaystyle \frac{dT_i}{d\varphi }}\delta (uu_i)`$ (15)
The delta functions (due to the presence of the branes) imply that $`\omega ^{}`$ and $`\varphi ^{}`$ are discontinuous at the branes and satisfy the matching conditions
$`3(\omega ^{}(u_i+0)\omega ^{}(u_i0))`$ $`=`$ $`\kappa ^2T_i|_{u_i}`$
$`Z|_{u_i}(\varphi ^{}(u_i+0)\varphi ^{}(u_i0))`$ $`=`$ $`\kappa ^2{\displaystyle \frac{dT_i}{d\varphi }}|_{u_i}`$ (16)
It should be noted that general covariance would imply that the scalar field equation should be satisfied when Einstein’s equations (the first two in the above set (3)) are satisfied. In the presence of the branes (which break the five dimensional general covariance) the consistency of the third with the first two implies a condition
$$(\varphi ^{}\frac{dT_i}{d\varphi })|_{u_i}=4(\omega ^{}T_i)|_{u_i}$$
(17)
where we may define $`\varphi ^{}(u_i)=\frac{1}{2}(\varphi ^{}(u_i+0)+\varphi ^{}(u_i0))`$ and similarly with $`\omega ^{}(u_i)`$. In fact this condition is the same as what one would get from requiring that the potential be continuous at $`u=u_i`$ and using the first equation of (3). However when $`\varphi ^{}(0),\omega ^{}(0)`$ are zero (as is the case if we impose a $`Z_2`$ symmetry under $`uu`$) (17) is trivially satisfied.
Let us first consider solutions with one brane located say at $`u=0`$. Also suppose that the bulk potential is of the form
$$V=\left(\frac{W}{\varphi }\right)^2\frac{4}{3}W^2.$$
(18)
where $`W=W(\varphi )`$ may be considered to be a sort of superpotential. This form for $`V`$ arises naturally in gauged supergravities and appears to be a necessary condition for the existence of a solution ,,. In this case the solutions for the warp factor and the scalar field can be obtained from ,,
$$3\omega ^{}=W(\varphi ),\varphi ^{}=\frac{dW}{d\varphi }$$
(19)
which can be solved by quadratures. Given these bulk solutions then the existence of a flat brane is guaranteed provided the matching condition is satisfied. But this is just a matter of choosing integration constants.
Let us discuss this further. We will impose a $`Z_2`$ symmetry as in ,. This might be a useful constraint in that the most likely string realization of the brane world scenario is probably a type II orientifold. Thus we impose
$$\omega (u)=\omega (u),\varphi (u)=\varphi (u).$$
(20)
The matching conditions (3) become (for the brane at $`u=0`$),
$`3\omega ^{}|_{u=0+}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\kappa ^2T_0(\varphi )|_{u=0+}`$
$`Z(\varphi )\varphi ^{}|_{u=0+}`$ $`=`$ $`+{\displaystyle \frac{1}{2}}\kappa ^2{\displaystyle \frac{dT_0}{d\varphi }}(|_{u=0+}`$ (21)
The two second order equations for $`\omega `$ and $`\varphi `$ would have two integration constants each. However the first equation of (3) is an energy integral with the total energy being zero. So the number of constants is reduced to three. Also a constant in $`\omega `$ is irrelevant since the equations of motion do not involve $`\omega `$ (this reflects the fact that such a constant can be absorbed in the rescaling of coordinates). Thus there really are only two constants (say $`\varphi (0)`$ and $`\omega ^{}(0)`$) that can be then chosen to satisfy the matching conditions. As explained in when the first order equations in terms of $`W`$ are being solved one would replace $`\omega ^{}(0)`$ by the integration constant coming from integrating (18). Thus with one brane a flat solution can be obtained without any fine tuning. Such a one brane solution we believe is unlikely to arise say from string theory since the brane typically carries some charge which would mean that the fifth dimension would have to be non-compact. This may however be a way of getting the scenario of the second paper of , but with the exponential potential for $`\varphi `$ that naturally arises in string theory, one gets a logarithmic behaviour for the warp factor $`\omega `$ ,<sup>7</sup><sup>7</sup>7In the last two references the singularity in such a metric is interpreted as a point where the space is to be cut off. However it is not entirely clear to us how this can arise from a microphysical theory such as string theory. rather than the linear behaviour required in . Later on we will come back to the scenario of in a situation where the modulus field has been integrated out from the low energy theory.
When there are two branes there is another pair of matching conditions to satisfy, but also there is another parameter namely the value $`u=R`$ at which the new brane is situated (so in the IIB example this would be size of the orbifolded fifth dimension). Then we have an additional pair of conditions,
$`3\omega ^{}|_{u=R}`$ $`=`$ $`+{\displaystyle \frac{1}{2}}\kappa ^2T_0(\varphi )|_{u=R}`$
$`Z(\varphi )\varphi ^{}|_{u=R}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\kappa ^2{\displaystyle \frac{dT_0}{d\varphi }}(|_{u=R}`$ (22)
From the energy constraint (the first equation of (3) we also have
$$Z^1\left(\frac{\kappa ^2}{2}\frac{dT_0}{d\varphi }\right)^2|_{u=0+}\frac{4}{3}\left(\frac{\kappa ^2}{2}T_0\right)^2|_{u=0+}=V|_{u=0+}$$
(23)
There is of course a similar equation at the point $`u=R`$ but this is not independent. Once we have a solution to the equations of motion and the matching conditions this will be automatically satisfied. In general the last equation will have a discrete set of solutions for $`\varphi (0)`$.
Thus there is one extra condition and to satisfy that requires a fine tuning, either of the brane tension or of the potential . However here too fine tuning can be avoided if we make at least one coupling constant in the potential dynamical, i.e. an integration constant. This is easily done if the bulk potential comes (at least partly) from the five dimensional analog of the $`F^2`$ term in (1) or (9). In this case, after solving the $`F`$ equation of motion as in the discussion in section I and substituting to get an effective action without $`F`$, one gets a potential for $`\varphi `$ which depends on the integration constant $`\mu `$ as in the discussion after (1). In the string theory context however (as we shall discuss in more detail later on) the brane is a source for the five form field and so $`\mu `$ is quantized. However in that case as we shall see later, if one considers (squashed) sphere compactifications one gets an additional adjustable parameter which is an integration constant of the ten dimensional theory. This is of course fine tuning from the from the point of view of the five dimensional theory, but not from the fundamental ten dimensional point of view.
There are several different cases one may consider.
a) Supersymmetry is unbroken both in bulk and on brane(s).
b) Supersymmetry is preserved in bulk and broken on the brane(s).
c) Supersymmetry broken in both bulk and brane(s).
d) Dilaton (and all other moduli) are fixed at the string scale
Let us discuss in turn the above cases.
a) Let us for example take a case which can come from type IIB orientifold constructions compactified to five dimensions (say on a torus or an orbifold). The low energy effective action contains a term $`U(\varphi )F_5^{}F_5`$, where $`\varphi `$ is the five dimensional dilaton. If one solves the equation of motion for the corresponding gauge field as in (10) then one effectively gets a potential $`V(\varphi )=\mu ^2U^1(\varphi )`$. Thus in the type IIB case, $`V=\mu ^2exp(\frac{5}{3}\varphi )`$.In this case $`T_i=\tau _ie^{\frac{5}{6}\varphi }`$ where $`\tau _i`$ is the (constant) brane charge. If we substitute this in (23) we see that in fact the $`\varphi (0)`$ dependence drops out and the equation is satisfied if $`\mu ^2=\frac{13}{9}(\frac{\kappa ^2}{2})^2\tau ^2`$. It should be noted that in this case even with one brane one needs the non-zero solution to the $`F`$ equation (i.e. (10). This is to be expected since as one crosses the brane the $`F`$ field must change by the number of branes times the charge on a brane and in the supersymmetric case this charge is related to the tension.<sup>8</sup><sup>8</sup>8This might be thought of as being similar to that discussed in and . In fact in the latter paper it is pointed out that this case does not exist in source free gauged supergravity. However our case is somewhat different in that we have explicit sources. So it is not completely clear to us that the arguments of apply. In the two brane case there is no determination of the distance between the branes as is to be expected.
b) This case is more interesting. Now supersymmetry is broken on the brane and so the tension need not be as in a). In this case one would expect (23) to determine $`\varphi (0)`$ and the matching conditions will determine the other two integration constants. In the first order formalism one of the integration constants will be the value of (say) $`W(\varphi (0))`$. If we work in the second order formalism after fixing $`\varphi (0)`$ as above the two constants to be determined by the the two matching conditions (3) are $`\varphi ^{}(0)`$ and $`\omega ^{}(0)`$). Thus one would indeed obtain (by choice of integration constants) a flat brane in 4D without fine tuning. When there is a second brane however, as we discussed earlier, there is one extra parameter (the distance $`R`$), but two more matching conditions to satisfy, and so we need to have a dynamical bulk cosmological constant.
c) In this case the bulk potential will also get renormalized but as far as the existence of flat brane solutions without fine tuning goes, there are no qualitatively new features compared to b).
d) This case we believe is quite interesting since it seems very likely that the moduli are fixed at (or close to) the string scale.<sup>9</sup><sup>9</sup>9For a discussion of this with references to earlier work see . This as we mentioned earlier would correspond to the original Randall-Sundrum scenario . This is possible in a situation in which stringy non-perturbative effects give a potential to all the moduli which should therefore be integrated out from the low energy effective action. In the absence of a string field theory, it is difficult to make precise statements . It is hard to see how the integration constant we want would arise if all moduli are fixed at the string scale since all we know how to deal with are low energy actions. Nevertheless perhaps one can make some educated guesses about this case too.
A possible string theoretic construction for the scenarios in cases b) and c) and possibly d) with one modulus field (the breathing mode) remaining in the action, may run as follows. Consider type IIB compactified to five dimensions on $`S^5`$ as in (section 2.4). We assume that the ten dimensional dilaton $`\phi `$ is fixed by string scale dynamics, but we keep the breathing mode that arises in the compactification, in the action. The relevant part of the low energy effective action is
$$S=d^5x\sqrt{G}(R(\varphi )^2)8m^2e^{8\alpha \varphi }+e^{\frac{16\alpha }{5}\varphi }R_5)$$
(24)
In getting the above from the 10 D action (or equations of motion) of IIB supergravity, the ansatz
$$ds_{10}^2=e^{2\alpha \varphi }ds_5^2+e^{2\beta \varphi }ds^2(S^5)$$
with $`\alpha =\frac{1}{4}\sqrt{\frac{5}{3}},\beta =\frac{3}{5}\alpha `$ has been made. $`R_5`$ is the Ricci scalar of the compact 5-manifold (in this case $`S^5`$) and is an unconstrained positive constant. The ansatz for the self-dual 5-form, is
$$F_5=4me^{8\alpha \varphi }ϵ_{(5)}+4mϵ_{(5)}(S^5)$$
(25)
where $`m`$ is an integration constant.
At this point let us show out how the quantization of $`m`$ mentioned earlier comes about if this scenario emerges from string theory. In this case the sources of the $`F_5`$ are $`D_3`$ branes (and/or orientifolds). The consistency of the the coupling of these objects leads to the standard Dirac quantization rule,
$$\tau _3_{S_5}F_5=2\pi n$$
where $`n`$ is an integer. Using the ansatz (25) for $`F_5`$ then gives $`\tau _34m\pi ^3r^5=2\pi n`$ where $`r`$ is the radius of $`S_5`$. This then gives us (after using the D3-brane tension formula $`\tau _3=2\pi M_s^4`$)
$$m=\frac{n}{4\pi ^3M_s^4r^5}=\frac{nM_s}{4\pi ^3\widehat{r}^5}$$
(26)
where we have introduced the dimensionless parameter $`\widehat{r}rM_s`$. Also the Ricci scalar of $`S_5`$ may then be written as $`R_5=\frac{20M_s^2}{\widehat{r}^2}`$. Then the potential in (24) may be rewritten
$`V`$ $`=`$ $`M_s^2\left[{\displaystyle \frac{n^2}{2\pi ^6\widehat{r}^{10}}}e^{8\alpha \varphi }{\displaystyle \frac{20}{\widehat{r}^2}}e^{\frac{16}{5}\alpha \varphi }\right]`$ (27)
$`=`$ $`M_s^2\widehat{r}^{\frac{10}{3}}\left[{\displaystyle \frac{n^2}{2\pi ^6}}e^{8\alpha \stackrel{~}{\varphi }}20e^{\frac{16}{5}\alpha \stackrel{~}{\varphi }}\right]`$
where in the last line of the above equation we put $`\alpha \varphi =\alpha \stackrel{~}{\varphi }+\frac{5}{3}\mathrm{ln}\widehat{r}`$.
The minimum of this potential is given by $`e^{\frac{24}{5}\alpha \stackrel{~}{\varphi }_0}=\frac{16\pi ^6}{n^2}`$. As discussed in the 5D action allows a $`AdS_5`$ solution with the cosmological constant being an integration constant, with the breathing mode being fixed at ($`\stackrel{~}{\varphi }_0`$). The potential at this point gives a cosmological constant $`\mathrm{\Lambda }=V(\stackrel{~}{\varphi }_0)=(2\pi )^4\frac{16}{n^2}M_s^2\widehat{r}^{\frac{10}{3}}`$. It should be noted that this five dimensional cosmological constant is dependent on the compactification parameter $`\widehat{r}`$ and can be adjusted even though there are no adjustable constants in the 10 dimensional action.
Let us now compactify one of the spatial dimensions on $`S^1/Z_2`$ as before. In the original string theory scenario we would have $`D3`$ branes and orientifolds distributed over the five sphere but we will just consider the effective theory in five dimensions with two branes sitting at the ends of the compact fifth dimension giving us the scenario we had earlier.
The main focus of our discussion is going to be on how to get a two -brane scenario without fine tuning of parameters in the 10 dimensional Lagrangian. This will require that we keep the breathing mode $`\stackrel{~}{\varphi }`$ as a dynamical field that is not sitting at the minimum of the potential. However before we do that let us see whether we can get a Randall-Sundrum type scenario with the scalar breathing mode integrated out. For this we need to assume that some string non-perturbative effects fix the breathing mode at some high scale so that in the effective low energy five dimensional theory it has been integrated out just like the ten dimensional dilaton. Then we would have a bulk cosmological constant which will get contributions from the string scale effects as well as from V. The latter will of course not necessarily be fixed at $`V(\stackrel{~}{\varphi }_0)`$ since it will be primarily determined by string effects however the important point is that such a contribution will be $`\widehat{r}`$ dependent and hence will be adjustable. Making the mild assumption that this effective cosmological constant is negative we put $`V=\frac{\mu ^2}{4}(\widehat{r})`$ ) in the first equation of (3). This then becomes $`12\omega ^2=\frac{\mu ^2}{4}`$ giving $`\omega ^{}=\pm \frac{\mu }{4\sqrt{3}}`$, so that $`\omega =\frac{\mu }{4\sqrt{3}}|u|`$. In the last equation we have used the $`Z_2`$ symmetry so as to obtain a warp factor that decays exponentially from the origin on both sides. Using the matching condition (3) then gives $`\kappa ^2T_0=\frac{3\mu }{2}`$. The point is that this condition can be satisfied without fine tuning of the tension since the bulk potential is a function of the adjustable compactification constant $`\widehat{r}`$. If the size of our compact direction is taken to infinity then we have the RS2 scenario. On the other hand if this dimension is finite then we need a second brane at $`u=R`$ necessarily its tension is negative $`T_1=\frac{3\mu }{2}`$. This of course is then a fine tuning.
There are several points that should be noted in this calculation.
* In the absence of the modulus field there is no flat one brane solution without fine tuning (as in ) or having a dynamical cosmological constant as in the above discussion. Indeed in the latter case there is (perhaps) a theory of confined gravity as in the second paper of but obtained now without fine tuning of the fundamental (ten-dimensional) Lagrangian.
* In the RS1 scenario the distance $`R`$ is now a free parameter (adjusted to a value that “explains” the gauge hierarchy in ) and is not fixed by the dyanmics. Indeed the scalar field was introduced in in order to stabilize the value of $`R`$. However this requires a tuning of a parameter in the potential in order to obtain the “right” value. So unless this value of the parameter in the potential has a natural explanation there is no particular advantage to this.
* In the two brane case the so-called visible brane (on which the standard model is supposed to live) has negative tension. Also since the dynamical bulk cosmological constant tracks the brane tension at the origin as it changes with RG scale the only way (without fine tuning) for a two brane solution to be viable is for the RG flow of the visible brane to be the same in magnitude though opposite in sign as on the other brane. It is not clear to us how to achieve this in a natural way. Thus this still requires fine tuning.
Let us now get back to the main point of this discussion. To discuss a two brane case (in cases b),d) above) with at least one modulus field dynamical we should look at a IIB orientifold. Now the fifth dimension is an interval $`S_1/Z_2`$ with 16 orientifold fixed planes at the fixed points $`u=0`$ and $`u=R`$. One also needs to introduce D-branes in order to cancel tadpoles<sup>10</sup><sup>10</sup>10Indeed such a model is T-dual to the type IA theory discussed in .. In the presence of D-branes and orientifold planes that are charged under the $`F_5`$ field we have (as in ) a discontinuity in $`m`$ by an amount equal to the brane charge/tension at the position of the brane. Imposing also the $`Z_2`$ symmetry the constant $`m`$ in (24) would be fixed in terms of the brane charge as in . This is of course just another way of verifying the quantization discussed earlier. When supersymmetry is broken however the brane tension would get renormalized so that the supersymmetric relation between tension and charge will be lost. Nevertheless the constant $`\widehat{r}`$ can adjust itself now to track the brane tension. In addition (assuming it is not fixed at the string scale by stringy effects as in the previous example) we have a modulus field $`\varphi `$ as in the discussion at the begining of this section, to supply an addtional integration constant so that one may have solutions with two flat branes.
Let us give some more details of this scenario. The effective five dimensional theory in the presence of branes is given by (after putting in also the value of the five D Newton constant)
$`S`$ $`=`$ $`2\pi ^4M_s^3\widehat{r}^5{\displaystyle d^5x\sqrt{g_5}\left(R\frac{1}{2}(\stackrel{~}{\varphi })^2+M_s^2\widehat{r}^{\frac{10}{3}}U(\stackrel{~}{\varphi })\right)}`$ (28)
$`+`$ $`{\displaystyle \frac{2\pi M_s^4}{\pi ^3}}\left\{{\displaystyle _{u=0}}d^4x\sqrt{g_4}T_0(\stackrel{~}{\varphi }){\displaystyle _{u=R}}d^4x\sqrt{g_4}T_1(\stackrel{~}{\varphi })\right\}.`$
In the above we have written the potential $`V=M_s^2R^{\frac{10}{3}}U(\stackrel{~}{\varphi })`$ where $`U`$ is independent of $`\widehat{r}`$ (see (27)). Let us first look at what might happen in the supersymmetric case. Here in analogy with the case discussed by Lucas et al for the Supersymmetry preserving compactification of M theory on $`\frac{S^1}{Z_2}\times CY_3`$ we expect the BPS equations
$$T_0(\stackrel{~}{\varphi })=T_1(\stackrel{~}{\varphi })=\pi ^3\widehat{r}^{5/3}W(\stackrel{~}{\varphi })$$
where
$$U(\stackrel{~}{\varphi })=\left(\frac{W}{\stackrel{~}{\varphi }}\right)^2\frac{2}{3}W^2$$
. (Note that apart from normalization this is the same $`W`$ as before). In fact these BPS conditions have been justified for this system very recently in \[bkv\],\[dls\]. In this case it is easily seen that the parameter $`\widehat{r}`$ drops out of the matching conditions which however can be satisfied for arbitrary inter-brane distance $`R`$ and integration constants. This is exactly like what happens in the M theory case investigated by Lukas et al. \[lukas\].
In the broken supersymmetric case however the situation is quite different. For instance we may imagine now is that the compactification is on a squashed sphere so that the supersymmetry is $`N=1`$ which is then broken by for instance by gaugino condensation effects on one or other brane. Now the tension will be renormalized from its BPS value so that one would expect
$$T_i(\stackrel{~}{\varphi })\pm \pi ^3\widehat{r}^{\frac{5}{3}}W(\stackrel{~}{\varphi })+ϵ\psi _i(\widehat{r},\stackrel{~}{\varphi })$$
Here the function $`\psi `$ is dependent on the details of the low energy field theories on the branes. Thus now $`\widehat{r}`$ will not drop out of the matching conditions and indeed we need to adjust the two integration constants $`\stackrel{~}{\varphi }(0),A^{}(0)`$ the inter-brane distance $`R`$ as well as the compactification parameter $`\widehat{r}`$ in order to get a flat four-dimensional brane world.
A detailed discussion of such models will be given in a forthcoming paper .
## 4 Conclusions
Let us first discuss the results of ,. From our discussion it should be clear that the reason that flat (one) brane solutions are obtained (without fine tuning) in these works is that integration constants have been chosen to ensure the existence of such solutions. Of course since these authors do not discuss two brane solutions they do not need the $`F_5`$ field or the Ricci non-flat compactification that we have introduced. However as we argued (and is indeed implied by the work of de Wolf et al ) one flat brane solution is obtained in the presence of a dynamical scalar field by choosing the integration constant $`\varphi (0)`$ appropriately. It does not depend on the particular form of the brane tension $`T(\varphi )`$ as seems to be implied in . Indeed this is just as well since the form of this function can change under renormalization effects on the brane. The fact that only a flat brane is allowed for a particular form of this function (see equation (14) of and section (3.2) of ) therefore is not a RG invariant statement. Quantum effects of the standard model in a background metric yields both a cosmological constant as well as curvature terms (as in our (7)). The latter will necessarily modify these arguments.
The main conclusion of the present work is that one can indeed obtain flat branes (and in particular zero cosmological constant in the brane containing the standard model) without fine tuning, but it involves a choice of integration constants/compactification parmeters. In this respect these theories have the same problem that bedevilled those of references , ,. It is useful to review this issue briefly. The point is to show that the particular integration constant(s) that leads to a zero cosmological constant gets chosen because it is the most probable one. To show this Hawking used a Euclidean quantum gravity argument according to which (see also section VIII of Weinberg’s classic review ) the probability for the occurence of a value $`\mu `$ for the integration constant was given by $`P(\mu )\mathrm{exp}(\mathrm{\Gamma }_E[\psi _c])`$ where $`\mathrm{\Gamma }_E`$ is the Euclidean quantum effective action (essentially our equation (9) Wick rotated to a Euclidean metric) and the $`\psi _c`$ are the values of all the fields evaluated at an extremum of $`\mathrm{\Gamma }`$. The Euclidean (effective) action for a $`D`$ dimensional theory after setting all other fields but the metric to their quantum ground state values as above would take the form, (setting the $`D`$ dimensional Planck mass equal to one) $`\mathrm{\Gamma }_E=\sqrt{G}(R\mathrm{\Lambda })`$. From the Einstein equation we have $`R=\frac{D\mathrm{\Lambda }}{(D2)}`$. Substituting this into the Euclidean action gives $`S_E=\frac{2V_D}{D2}\mathrm{\Lambda }`$ where $`V_D`$ is the volume of Euclidean D space. If $`\mathrm{\Lambda }`$ is positive then the space is $`S_D`$ and its volume is $`V_D=\frac{a^2}{\mathrm{\Lambda }^2}`$ so that the action becomes $`S_E=\frac{2V}{(D2)\mathrm{\Lambda }}`$. Thus the probability distribution becomes
$$P(\mu )e^{\mathrm{\Gamma }_E[\psi _c]}=e^{+\frac{2V}{(D2)\mathrm{\Lambda }}}.$$
(29)
This would imply that the probability was peaked at $`\mathrm{\Lambda }0+`$. <sup>11</sup><sup>11</sup>11It was pointed out by Duff that if one substitutes the solution for $`F`$ into Einstein’s equation and then infers the effective action from which it comes, then one finds in fact that the (Euclidean) action is positive near $`\mathrm{\Lambda }=0`$ so that this value is actually disfavoured! Nevertheless it was claimed in that the correct action has a boundary term that does not affect the equation of motion and its inclusion will reestablish Hawking’s result. I wish to thank R. Bousso for pointing out this reference to me. On the other hand as pointed out by M. Duff (private communication) Hawking’s Euclidean space no-boundary action should not of course have a boundary contribution! The situation therefore remains murky.
In our case it is not clear whether an analog of Hawking’s argument would work<sup>12</sup><sup>12</sup>12Hawking’s argument may work also in the case of unimodular gravity .. However one would think that one should apply the argument to the five (or ten?) dimensional theory since that is the action one is starting from. But the integration constants must get chosen so that it is the four dimensional theory that has to have zero cosmological constant. At this point it is not clear to us whether a version of this argument can be used to justify the choice of integration constants.
Note added: While this paper was being prepared for submission, a paper which (inter alia) makes comments related to ours about the one brane case of ,, appeared as an e-print .
## 5 Acknowledgements
I wish to thank Alex Flournoy and Nicos Irges for discussions and Shamit Kachru for correspondence on the work of ,. I’m also grateful to Jack Ng and Paul Townsend for setting me straight on the historical record, and Renata Kallosh for correspondence on . This work is partially supported by the Department of Energy contract No. DE-FG02-91-ER-40672. |
warning/0002/math0002108.html | ar5iv | text | # Rolle’s theorem is either false or trivial in infinite-dimensional Banach spaces
## 1. Introduction and main results
Rolle’s theorem in finite-dimensional spaces states that, for every bounded open subset $`U`$ of $`^n`$ and for every continuous function $`f:\overline{U}`$ such that $`f`$ is differentiable in $`U`$ and constant on the boundary $`U`$, there exists a point $`xU`$ such that $`f^{}(x)=0`$. Unfortunately, Rolle’s theorem does not remain valid in infinite dimensions. It was S. A. Shkarin that first showed the failure of Rolle’s theorem in superreflexive infinite-dimensional spaces and in non-reflexive spaces which have smooth norms. The class of spaces for which Rolle’s theorem fails was substantially enlarged in , where it was also shown that an approximate version of Rolle’s theorem remains nevertheless true in all Banach spaces. In fact, as a consequence of the existence of diffeomorphisms deleting points in infinite-dimensional spaces (see ), it is easy to see that Rolle’s theorem fails in all infinite-dimensional Banach spaces which have smooth norms .
However, none of these results allows to characterize the spaces for which Rolle’s theorem fails since, as shown by R. Haydon , there are Banach spaces with smooth bump functions which do not possess any equivalent smooth norms. Of course, Rolle’s theorem is trivially true in the Banach spaces which do not have any smooth bumps (if $`X`$ is such a space then every function on $`X`$ satisfying the hypothesis of Rolle’s theorem must be a constant). Thus, in many infinite-dimensional Banach spaces, Rolle’s theorem either fails or is trivial, depending on the smoothness properties of the spaces considered. In this setting, it does not seem too risky to conjecture, as it was done in , that Rolle’s theorem should fail in an infinite-dimensional Banach space if and only if our space has a $`C^1`$ smooth bump function. In this paper we will prove this conjecture to be right, thus providing an interesting new characterization of smoothness in Banach spaces.
Our main result is the following
###### Theorem 1.1.
Let $`X`$ be an infinite-dimensional Banach space which has a $`C^p`$ smooth (Lipschitz) bump function. Then there exists another $`C^p`$ smooth (Lipschitz) bump function $`f:X[0,1]`$ with the property that $`f^{}(x)0`$ for every $`x\text{int}(\text{supp}f)`$.
Here $`\text{supp}f`$ denotes the support of $`f`$, that is, $`\text{supp}f=\overline{\{xX:f(x)0\}}`$. Let us recall that $`b:X`$ is said to be a bump function on $`X`$ provided $`b`$ is not constantly zero and $`b`$ has a bounded support.
From this result it is easily deduced the following
###### Corollary 1.2.
Let $`X`$ be an infinite-dimensional Banach space. The following statements are equivalent.
* $`X`$ has a $`C^p`$ smooth (and Lipschitz) bump function.
* There exist a bounded contractible open subset $`U`$ of $`X`$ and a continuous function $`f:\overline{U}`$ such that $`f`$ is $`C^p`$ smooth (and Lipschitz) in $`U`$, $`f=0`$ on $`U`$, and yet $`f^{}(x)0`$ for all $`xU`$, that is, Rolle’s theorem fails in $`X`$.
* There exist a $`C^p`$ smooth (and Lipschitz) function $`f:X[0,1]`$ and a bounded contractible open subset $`U`$ of $`X`$ such that $`f=0`$ precisely on $`XU`$ and yet $`f^{}(x)0`$ for all $`xU`$.
Just in order to complete the picture of Rolle’s theorem in infinite-dimensional Banach spaces let us quote the two positive results from on approximate and subdifferential substitutes of Rolle’s theorem, which guarantee the existence of arbitrarily small derivatives (instead of vanishing ones) for every function sayisfying (in an approximate manner) the conditions of the classic Rolle’s theorem.
###### Theorem 1.3 (Azagra–Gómez–Jaramillo).
Let $`U`$ be a bounded connected open subset of a Banach space $`X`$. Let $`f:\overline{U}`$ be a bounded continuous function which is (Gâteaux) differentiable in $`U`$. Let $`R>0`$ and $`x_0U`$ be such that $`\text{dist}(x_0,U)=R`$. Suppose that $`f(U)[\epsilon ,\epsilon ]`$ for some $`\epsilon >0`$. Then there exists some $`x_\epsilon U`$ such that $`f^{}(x_\epsilon )\frac{\epsilon }{R}`$.
###### Theorem 1.4 (Azagra–Deville).
Let $`U`$ be a bounded connected open subset of a Banach space $`X`$ which has a $`C^1`$ smooth Lipschitz bump function. Let $`f:\overline{U}`$ be a bounded continuous function, and let $`R>0`$ and $`x_0U`$ be such that $`\text{dist}(x_0,U)=R`$. Suppose that $`f(U)[\epsilon ,\epsilon ]`$ for some $`\epsilon >0`$. Then
$$inf\{p:pD^{}f(x)D^+f(x),xU\}\frac{2\epsilon }{R}.$$
(Here $`D^{}f(x)`$ and $`D^+f(x)`$ denote the subdifferential and superdifferential sets of $`f`$ at $`x`$, respectively; see , p. 339, for the definitions).
The “twisted tube” method that we develop in section 2 in order to prove theorem 1.1 is interesting in itself and, with little more work, provides a useful characterization of $`C^p`$ smoothness in infinite-dimensional Banach spaces related to the existence of a certain kind of deleting diffeomorphisms. Namely, we have the following
###### Theorem 1.5.
Let $`X`$ be an infinite-dimensional Banach space. The following assertions are equivalent.
* $`X`$ has a $`C^p`$ smooth bump function.
* There exists a nonempty contractible closed subset $`D`$ of the unit ball $`B_X`$ and a $`C^p`$ diffeomorphism $`f:XXD`$ so that $`f`$ restricts to the identity outside $`B_X`$.
This result yields the following corollaries.
First, the celebrated Brouwer’s fixed point theorem fails even for smooth self-mappings of balls or starlike bodies in all infinite-dimensional Banach spaces. Let us recall that Brouwer’s theorem states that every continuous self-map of the unit ball of a finite-dimensional normed space admits a fixed point. This is the same as saying that there is no continuous retraction from the unit ball onto the unit sphere, or that the unit sphere is not contractible (the identity map on the sphere is not homotopic to a constant map). In infinite dimensions the situation is completely different and Brouwer’s theorem is no longer true (see . Theorem 1.5 yields a trivial proof that Brouwer’s theorem is false in infinite dimensions even for smooth self-mappings of balls or starlike bodies; this is a particular case (the non-Lipschitz one) of the main result in .
Second, we deduce from the above characterization that the support of the bump functions which violate Rolle’s theorem can always be assumed to be a smooth starlike body. This is all shown in section 3.
In section 2 we give the proofs of theorems 1.1 and 1.5. A much simpler proof of theorem 1.1 for the non-Lipschitz case is included in this section too.
Finally, in section 4 we study the structure of the set of gradients of bump functions in the Hilbert space $`\mathrm{}_2`$, and as a consequence of the failure of Rolle’s theorem in infinite dimensions we get the following result. The usual norm of the Hilbert space $`\mathrm{}_2`$ can be uniformly approximated by $`C^1`$ smooth Lipschitz functions $`\psi `$ so that the cones generated by the sets of derivatives $`\psi ^{}(\mathrm{}_2)`$ have empty interior. This implies that there are $`C^1`$ smooth Lipschitz bumps in $`\mathrm{}_2`$ so that the cones generated by their sets of gradients have empty interior.
## 2. The proofs
The idea behind the proof of theorem 1.1 is as simple as this. First we build a twisted tube $`T`$ of infinite length in the interior of the unit ball $`B_X`$, with a beginning but with no end. This twisted tube can be thought of as directed by an ever-winding infinite path $`p`$ that gets lost in the infinitely many dimensions of our space $`X`$. In technical words, one can construct a diffeomorphism $`\pi `$ between a straight (unbounded) half-cilynder $`C`$ and a twisted (bounded) tube $`T`$ contained in $`B_X`$. The tube $`T`$ is going to be the support of a smooth bump function $`f`$ that does not satisfy Rolle’s theorem. In order to define such a function $`f`$ we only have to make it strictly increase in the direction which is tangent to the leading path $`p`$ at each point of the tube $`T`$. The graph of $`f`$ would thus represent an ever-ascending stairway built upon our twisted tube, with a beginning but no end.
The spirit of the proof that (1) implies (2) in theorem 1.5 is not very different. We will make use of the diffeomorphism $`\pi `$ between a straight (unbounded) half-cilynder $`C`$ and a bounded twisted tube $`T`$ contained in $`B_X`$. If we consider a straight closed half-cilinder $`C^{}`$ contained in the interior of $`C`$ and directed by the same line as $`C`$, it is elementary that there is a diffeomorphism $`g:XXC^{}`$ so that $`g`$ restricts to the identity outside $`C`$. In fact this is true even in the plane. Now, by composing this diffeomorphism $`g`$ with the diffeomorphisms $`\pi `$ and $`\pi ^1`$ that give us an appropriate coordinate system in the twisted tube $`T=\pi (C)`$, we get a diffeomorphism $`f:XXT^{}`$, where $`T^{}=\pi (C^{})`$ is a smaller closed twisted tube inside $`T`$, and $`f`$ restricts to the identity outside the unit ball. The precise definition of $`f`$ would be $`f(x)=\pi (g(\pi ^1(x)))`$ if $`xT`$, and $`f(x)=x`$ if $`xXT`$. If we take $`D=T^{}`$ we are done.
In the rest of his section we will be involved in the task of formalizing these ideas.
The following theorem guarantees the existence of bounded infinite twisted tubes in all infinite-dimensional Banach spaces.
###### Theorem 2.1.
There are universal constants $`M>0`$ (large) and $`\epsilon >0`$ (small) such that, for every infinite-dimensional Banach space $`X`$, if we consider the decomposition $`X=H[z]`$ (where $`H=\text{Ker}z^{}`$ for some $`z^{}X^{}`$ with $`z^{}(z)=z^{}=z=1`$) and the open half-cilynder $`C`$ of diameter $`\epsilon `$, directed by $`z`$, and with base on $`H`$, $`C=\{x+tzX:x<\epsilon ,t>0\}`$, then there exists an injection $`\pi :CB_X`$ which is a $`C^{\mathrm{}}`$ diffeomorphism onto its image. The image $`T=\pi (C)`$ is thus a bounded open set which we will call a bounded open infinitely twisted tube in $`X`$. Moreover, the derivatives of the mappings $`\pi :CT`$ and $`\pi ^1:TC`$ are both uniformly bounded by $`M`$.
Assume for a while that theorem 2.1 is already established and let us explain how theorems 1.1 and 1.5 can be deduced.
Proof of theorem 1.1.
Consider the diffeomorphism $`\pi :CTB_X`$ from theorem 2.1. Take a $`C^p`$ smooth (Lipschitz) non-negative bump function $`\phi `$ on $`H`$ so that the support of $`\phi `$ is contained in the base of $`C`$, that is, $`\phi (x)=0`$ whenever $`x\frac{\epsilon }{2}`$, for instance. Pick a $`C^{\mathrm{}}`$ smooth real function $`\mu :[0,1]`$ such that $`\mu (t)=0`$ for $`t1`$, $`0<\mu (t)<1`$ for $`t>1`$ and $`0<\mu ^{}(t)<1`$ for all $`t>1`$. Then define $`g:X=H[z]`$ by
$$g(x,t)=\phi (x)\mu (t).$$
It is plain that $`g`$ is a $`C^p`$ smooth (Lipschitz) function such that $`g^{}(x,t)0`$ for every $`x\text{int}(\text{supp}f)`$, that is, for every $`x`$ such that $`g(x,t)0`$ (take into account that the interior of the support of $`g`$ coincides in this case with the open support of $`g`$, that is the set of points at which $`g`$ does not vanish). Indeed,
$$g^{}(x,t)(0,1)=\frac{g}{t}(x,t)=\phi (x)\mu ^{}(t)$$
and therefore $`g^{}(x,t)(0,1)=0`$ if and only if $`\phi (x)=0`$ or $`\mu ^{}(t)=0`$, which happens if and only if $`\phi (x)=0`$ or $`\mu (t)=0`$, that is to say, $`g(x,t)=0`$. Now let us define $`f:X`$ by
$$f(y)=\{\begin{array}{cc}& g(\pi ^1(y))\text{if }yT\text{;}\hfill \\ & 0\text{if }yT\hfill \end{array}$$
It is clear that $`f`$ is a well defined $`C^p`$ smooth (Lipschitz) function, and $`\text{supp}(f)=\pi (\text{supp}(g))T`$, from which it follows that $`f`$ has a bounded support. We claim that $`f^{}(y)0`$ whenever $`y\text{int}(\text{supp}f)`$, that is, $`f`$ does not satisfy Rolle’s theorem. Indeed, if $`y\text{int}(\text{supp}f)`$ then $`\pi ^1(y)=(x,t)\text{int}(\text{supp}g)`$ and therefore $`g^{}(x,t)(0,1)0`$. But then
$$f^{}(y)=g^{}(x,t)D\pi ^1(y)0,$$
because $`D\pi ^1(y)`$ is a linear isomorphism. This concludes the proof of theorem 1.1.
Now we will turn our attention to the proof of theorem 1.5. Before proceeding with the proof, let us fix some standard terminology and notation used throughout this section and the following one. A closed subset $`A`$ of a Banach space $`X`$ is said to be a starlike body provided $`A`$ has a non-empty interior and there exists a point $`x_0\text{int}A`$ such that each ray emanating from $`x_0`$ meets the boundary of $`A`$ at most once. In this case we will say that $`A`$ is starlike with respect to $`x_0`$. When dealing with starlike bodies, we can always assume that they are starlike with respect to the origin (up to a suitable translation), and we will do so unless otherwise stated.
For a starlike body $`A`$, the characteristic cone of $`A`$ is defined as
$$ccA=\{xXrxA\text{for all}r>0\},$$
and the Minkowski functional of $`A`$ as
$$q_A(x)=inf\{\lambda >0\frac{1}{\lambda }xA\}$$
for all $`xX`$. It is easily seen that for every starlike body $`A`$ its Minkowski functional $`q_A`$ is a continuous function which satisfies $`q_A(rx)=rq_A(x)`$ for every $`r0`$ and $`q_A^1(0)=ccA`$. Moreover, $`A=\{xXq_A(x)1\}`$, and $`A=\{xXq_A(x)=1\}`$, where $`A`$ stands for the boundary of $`A`$. Conversely, if $`\psi :X[0,\mathrm{})`$ is continuous and satisfies $`\psi (\lambda x)=\lambda \psi (x)`$ for all $`\lambda 0`$, then $`A_\psi =\{xX\psi (x)1\}`$ is a starlike body. Convex bodies (that is, closed convex sets with nonempty interior) are an important kind of starlike bodies. We will say that $`A`$ is a $`C^p`$ smooth (Lipschitz) starlike body provided its Minkowski functional $`q_A`$ is $`C^p`$ smooth (and Lipschitz) on the set $`Xq_A^1(0)`$.
It is worth noting that for every Banach space $`(X,.)`$ with a $`C^p`$ smooth (Lipschitz) bump function there exist a functional $`\psi `$ and constants $`a,b>0`$ such that $`\psi `$ is $`C^p`$ smooth (Lipschitz) away from the origin, $`\psi (\lambda x)=|\lambda |\psi (x)`$ for every $`xX`$ and $`\lambda `$, and $`ax\psi (x)bx`$ for every $`xX`$ (see , proposition II.5.1). The level sets of this function are precisely the boundaries of the smooth bounded starlike bodies $`A_c=\{xX\psi (x)c\}`$, $`c`$. This shows in particular that every Banach space having a $`C^p`$ smooth (Lipschitz) bump function has a $`C^p`$ smooth (Lipschitz) bounded starlike body as well. The converse is obviously true too.
Proof of theorem 1.5.
First of all let us choose a number $`\epsilon >0`$, a cilynder $`C`$, a bounded twisted tube $`T`$, and a diffeomorphism $`\pi :CT`$ from theorem 2.1.
Let $`B`$ be a $`C^{\mathrm{}}`$ smooth convex body in the plane $`^2`$ whose boundary contains the set
$$\{(s,t):t=1,|s|\frac{\epsilon }{4}\}\{(s,t):|s|=\frac{\epsilon }{2},t1+\frac{\epsilon }{4}\},$$
and let $`q_B`$ be the Minkowski functional of $`B`$. Define $`B^{}=\frac{1}{2}B=\{(s,t):q_B(s,t)\frac{1}{2}\}`$. Let $`\theta :(\frac{1}{2},\mathrm{})[0,\mathrm{})`$ be a $`C^{\mathrm{}}`$ smooth real function so that $`\theta ^{}(t)<0`$ for $`\frac{1}{2}<t<1`$, $`\theta (t)=0`$ for $`t1`$, and $`lim_{t\frac{1}{2}^+}\theta (t)=+\mathrm{}`$. Now define $`\phi :^2B^{}^2`$ by
$$\phi (s,t)=(\phi _1(s,t),\phi _2(s,t))=(s,t+\theta (q_B(s,t))).$$
It is elementary to check that $`\phi `$ is a $`C^{\mathrm{}}`$ diffeomorphism from $`^2B^{}`$ onto $`^2`$ so that $`\phi `$ restricts to the identity outside the band $`B`$.
Next, recall that since $`X`$ has a $`C^p`$ smooth bump then it has a $`C^p`$ bounded starlike body $`A`$ as well. If $`X=H[z]`$, take $`W=AH`$, which is a $`C^p`$ bounded starlike body in $`H`$, and denote by $`q_W`$ its Minkowski functional. We can assume that $`WB(0,1)`$, that is, $`xq_W(x)`$ for all $`xH`$. Let us define
$$\psi (x,t)=q_B(q_W(x),t)$$
for all $`(x,t)X`$. It is clear that $`\psi `$ is a continuous function which is positive-homogeneous and $`C^p`$ smooth away from the half-line $`L=\{(x,t)X:x=0,t0\}`$. Then the sets
$$U=\{(x,t)X:\psi (x,t)1\},U^{}=\{(x,t)X:\psi (x,t)\frac{1}{2}\}$$
are cilyndrical $`C^p`$ starlike bodies whose characteristic cones are the half-line $`L`$. If we define
$$h(x,t)=(x,\phi _2^1(q_W(x),t))$$
for $`(x,t)X`$, it is not difficult to realize that $`h`$ is a $`C^p`$ diffeomorphism from $`X`$ onto $`XU^{}`$ so that $`h`$ restricts to the identity outside $`U`$. The inverse of $`h`$ is given by
$$h^1(x,t)=(x,t+\theta (\psi (x,t))).$$
Now consider the cilyndrical bodies $`V:=(0,2)+U`$ and $`V^{}:=(0,2)+U^{}`$, and put $`g(x,t)=h(x,t2)`$. Then $`g:XXV^{}`$ is a $`C^p`$ diffeomorphism such that $`g`$ is the identity outside $`V`$. Note that, since $`WB(0,1)`$, we have that $`V^{}VC=\{(x,t)X:x<\epsilon ,t>0\}`$. Let us define
$$f(x)=\{\begin{array}{cc}& \pi (g(\pi ^1(x)))\text{if }xT\text{;}\hfill \\ & x\text{otherwise.}\hfill \end{array}$$
It is then clear that $`f`$ is a $`C^p`$ diffeomorphism from $`X`$ onto $`XT^{}`$, where $`T^{}=\pi (V^{})`$ is a smaller closed twisted tube inside $`\pi (V)T`$, and $`f`$ restricts to the identity outside the larger tube $`\pi (V)T`$, which is contained in $`B_X`$. This completes the proof that (1) implies (2).
Conversely, if there is such an $`f`$ as in (2), we can assume that $`f(0)0`$ and take $`TX^{}`$ so that $`T(f(0))0`$; then the function $`b:X`$ defined by $`b(x)=T(xf(x))`$ is a $`C^p`$ smooth bump on $`X`$.
Now we proceed with the proof of theorem 2.1. We will make use of the following lemma, which guarantees the existence of an appropiate path of linear isomorphisms. Here $`\text{Isom}(X)`$ stands for the set of linear isomorphisms of $`X`$, which is regarded as a subset of $`(X,X)`$, the linear continuous mappings of $`X`$ into $`X`$.
###### Lemma 2.2.
There is a universal constant $`K>0`$ such that for every infinite-dimensional Banach space $`X`$ there are paths $`\beta :[0,\mathrm{})\text{Isom}(X)`$ and $`p:[0,\mathrm{})X`$ with the following properties:
* Both $`\beta `$ and $`p`$ are $`C^{\mathrm{}}`$ smooth, as well as the path of inverse isomorphisms $`\beta ^1:[0,\mathrm{})\text{Isom}(X)`$, $`\beta ^1(t)=[\beta (t)]^1`$.
* $`1\beta (t)K`$ and $`1\beta ^1(t)K`$ for all $`t[0,\mathrm{})`$.
* $`sup_{t0}\beta ^{}(t)K`$ and $`sup_{t0}(\beta ^1)^{}(t)K`$.
* There exists a certain $`vX`$, with $`1v\frac{1}{K}`$, such that $`p^{}(t)=\beta (t)(v)`$ for all $`t0`$.
* For every $`t,s[0,\mathrm{})`$ we have that $`p(t)p(s)\frac{1}{K}\mathrm{min}\{1,|ts|\}`$.
###### Proof.
Let $`(x_n)_{n=0}^{\mathrm{}}`$ be a normalized basic sequence in $`X`$ with biorthogonal functionals $`(x_n^{})_{n=0}^{\mathrm{}}X^{}`$ (that is, $`x_n^{}(x_k)=\delta _{n,k}=1`$ if $`n=k`$, and $`0`$ otherwise) satisfying $`x_n^{}3`$ (one can always take such sequences, see , p. 93, or , p. 39). For $`n1`$ set $`v_n=x_nx_{n1}`$. Let $`\theta :`$ be a $`C^{\mathrm{}}`$ function with the following properties:
* $`\theta (t)=0`$ whenever $`t\frac{1}{2}`$ or $`t1`$;
* $`\theta (t)=1`$ for $`t[0,\frac{1}{2}]`$;
* $`\theta ^{}(t)>0`$ for $`t(\frac{1}{2},0)`$;
* $`\theta (t)=1\theta (t1)`$ for $`t[\frac{1}{2},1]`$;
* $`sup_t|\theta ^{}(t)|4`$.
For $`n1`$ let us define $`\theta _n:`$ by $`\theta _n(t)=\theta (tn+1)`$. It is clear that the functions $`\theta _n`$ are all $`C^{\mathrm{}}`$ smooth and have Lipschitz constant less than or equal to $`4`$, $`\theta _n=0`$ on $`(\mathrm{},n1\frac{1}{2}][n,\mathrm{})`$, $`\theta _n=1`$ on $`[n1,n\frac{1}{2}]`$, and $`\theta _n(t)=1\theta _{n+1}(t)`$ for all $`t[n1,n+\frac{1}{2}]`$.
Our path $`\beta `$ of linear isomorphisms is going to be of the form
$$\beta (t)=\underset{n=1}{\overset{\mathrm{}}{}}\theta _n(t)S_n,$$
where each $`S_n\text{Isom}(X)`$ takes the vector $`v_1`$ into $`v_n`$ and for every $`\lambda [0,1]`$ the mapping $`L_{n,\lambda }=(1\lambda )S_n+\lambda S_{n+1}`$ is still a linear isomorphism and, moreover, the families of isomorphisms $`\{L_{n,\lambda }\}_{n,\lambda [0,1]}`$ and $`\{L_{n,\lambda }^1\}_{n,\lambda [0,1]}`$ are uniformly bounded. Let us define the isomorphisms $`S_n`$. They are going to be of the form
$$S_n(x)=x+f_n(x)(v_nv_1),$$
where $`f_nX^{}`$ satisfies $`f_n(v_1)=1=f_n(v_n)`$, and $`f_n18`$ (the exact definition of $`f_n`$ will be given later). Their inverses $`S_n^1`$ will be
$$S_n^1(y)=yf_n(y)(v_nv_1).$$
We want the linear mappings $`L_{n,\lambda }=(1\lambda )S_n+\lambda S_{n+1}`$ to be linear isomorphisms. We have
$`(1)`$
$$y=L_{n,\lambda }(x)=x+(1\lambda )f_n(x)(v_nv_1)+\lambda f_{n+1}(x)(v_{n+1}v_1),$$
from which
$`(2)`$
$$x=y[(1\lambda )f_n(x)(v_nv_1)+\lambda f_{n+1}(x)(v_{n+1}v_1)],$$
and we need to write $`f_n(x)`$ and $`f_{n+1}(x)`$ as linear functions of $`y`$. If we apply the functionals $`f_n`$ and $`f_{n+1}`$ successively to equation (1), we denote $`A_n=f_n(x)`$, $`B_n=f_{n+1}(x)`$, $`C_n=f_n(y)`$, $`D_n=f_{n+1}(y)`$, and we take into account that $`1=f_n(v_1)=f_n(v_n)=f_{n+1}(v_1)`$, then we obtain the system
$`(3)`$
$$\{\begin{array}{cc}& A_n+\lambda [f_n(v_{n+1})1]B_n=C_n\hfill \\ & (1\lambda )[f_{n+1}(v_n)1]A_n+B_n=D_n,\hfill \end{array}$$
which we want to have a unique solution for $`A_n,B_n`$. The determinant of this system is
$$\mathrm{\Delta }_{n,\lambda }=1\lambda (1\lambda )[f_{n+1}(v_n)1][f_n(v_{n+1})1],$$
and we want $`\mathrm{\Delta }_{n,\lambda }`$ to be bounded below by a strictly positive number, and this bound has to be uniform in $`n,\lambda `$. For $`n3`$ this can easily be done by setting
$$f_n=x_1^{}x_{n1}^{}$$
(so that $`f_n(v_n)=1=f_n(v_1)`$, $`f_n(v_{n+1})=0`$, $`f_{n+1}(v_n)=1`$, and therefore $`\mathrm{\Delta }_{n,\lambda }=(1\lambda )^2+\lambda ^2\frac{1}{2}`$ for all $`\lambda [0,1]`$). For $`n=1,2`$, put
$$f_2=x_1^{}+2x_2^{}+\frac{7}{3}x_3^{},\text{and}f_1=x_1^{};$$
then $`f_2(v_3)=\frac{1}{3}`$, $`f_2(v_2)=1`$, $`f_2(v_1)=1`$, $`f_3(v_2)=2`$, $`f_1(v_2)=1`$, $`f_1(v_1)=1`$, and everything is fine (indeed, $`\mathrm{\Delta }_{1,\lambda }=1`$ and $`\mathrm{\Delta }_{2,\lambda }=(1\lambda )^2+\lambda ^2\frac{1}{2}`$ for all $`\lambda [0,1]`$).
Therefore, with these definitions, the linear system (3) has a unique solution for $`A_n,B_n`$, which can be easily calculated and estimated by Cramer’s rule, of the form
$$\begin{array}{cc}& A_n(y)=\frac{1}{\mathrm{\Delta }_{n,\lambda }}\left(f_n(y)\lambda [f_n(v_{n+1})1]f_{n+1}(y)\right)\hfill \\ & B_n(y)=\frac{1}{\mathrm{\Delta }_{n,\lambda }}\left(f_{n+1}(y)(1\lambda )[f_{n+1}(v_n)1]f_n(y)\right).\hfill \end{array}$$
The linear forms $`yA_n(y)`$, $`yB_n(y)`$ satisfy that $`A_n144B_n`$ for all $`n`$, as is easily checked. Now, by substituting $`f_n(x)=A_n(y)`$ and $`f_{n+1}(x)=B_n(y)`$ in (2) we get the expression for the inverse of $`L_{n,\lambda }`$, that is,
$`(4)`$
$$x=L_{n,\lambda }^1(y)=y[(1\lambda )A_n(y)(v_nv_1)+\lambda B_n(y)(v_{n+1}v_1)].$$
By taking into account that $`A_n144B_n`$, $`f_n18`$ and $`v_nv_14`$ for all $`n`$, one can estimate that $`1L_{n,\lambda }73`$ and $`1L_{n,\lambda }^1577`$ for all $`n,\lambda [0,1]`$.
So let us define $`\beta :[0,\mathrm{})\text{Isom}(X)`$ by
$`(5)`$
$$\beta (t)=\underset{n=1}{\overset{\mathrm{}}{}}\theta _n(t)S_n.$$
This path is well defined because the sum is locally finite; in fact, from the definition of $`\theta _n`$ it is clear that, for a given $`t_0[0,\mathrm{})`$ there exist some $`\delta >0`$ and $`N=N(t_0)`$ such that $`\beta (t)=\theta _N(t)S_N+\theta _{N+1}(t)S_{N+1}`$ for all $`t(t_0\delta ,t_0+\delta )`$, that is, $`\beta `$ is locally of the form $`\beta (t)=L_{n,\lambda (t)}`$, where $`\lambda (t)=\theta _n(t)`$. This implies that the $`\beta (t)`$ are really linear isomorphisms and that the path is $`C^{\mathrm{}}`$ smooth.
On the other hand, the path $`\beta ^1(t)=[\beta (t)]^1\text{Isom}(X)`$ is $`C^{\mathrm{}}`$ smooth as well, because it is the composition of our path $`\beta `$ with the mapping $`\phi :\text{Isom}(X)\text{Isom}(X)`$, $`\phi (U)=U^1`$, which is $`C^{\mathrm{}}`$ smooth and whose derivative is given by $`\phi ^{}(U)(S)=U^1SU^1`$ for every $`S(X,X)`$ (see , theorem 5.4.3). This proves condition (1) of the lemma.
Next, by bearing in mind the local expression of $`\beta `$ and the above estimations for $`L_{n,\lambda }`$ and $`L_{n,\lambda }^1`$, we deduce that
$$1\beta (t)R\beta ^1(t)1$$
for all $`t[0,\mathrm{})`$, where $`R577`$ will be fixed later. This shows condition (2). Now, if $`t_0[0,\mathrm{})`$ and we write $`\beta (t)=\theta _N(t)S_N+\theta _{N+1}(t)S_{N+1}`$ for $`t(t_0\delta ,t_0+\delta )`$ as above, then it is clear that $`\beta ^{}`$ is locally of the form
$$\beta ^{}(t)=\theta _N^{}(t)S_N+\theta _{N+1}^{}(t)S_{N+1}$$
and therefore
$$\beta ^{}(t)|\theta _N^{}(t)|S_N+|\theta _{N+1}^{}(t)|S_{N+1}4(73+73)=584,$$
from which we get $`sup_{t0}\beta ^{}(t)584R`$. Moreover, we have
$$(\beta ^1)^{}(t)=(\beta (t))^1\beta ^{}(t)(\beta (t))^1$$
and therefore
$$(\beta ^1)^{}(t)\beta (t)^1^2\beta ^{}(t)(577)^2584,$$
from which $`sup_{t0}(\beta ^1)^{}(t)R`$ and condition (3) is satisfied as well provided we fix $`R=(577)^2584`$.
Now let us define the path $`p:[0,\mathrm{})X`$ by
$$p(t)=_{\mathrm{}}^t\beta (s)(v_1)𝑑s=_{\mathrm{}}^t\left(\underset{n=1}{\overset{\mathrm{}}{}}\theta _n(s)S_n(v_1)\right)𝑑s.$$
It is clear that $`p`$ is a $`C^{\mathrm{}}`$ smooth path in $`X`$, and $`p^{}(t)=\beta (t)(v_1)`$ for all $`t0`$ (from which it follows that $`p`$ is Lipschitz). Let us see that $`p`$ is bounded. For a given $`t>0`$ there exists $`N=N(t)`$ so that $`N1\frac{1}{2}tN\frac{1}{2}`$ and therefore, taking into account the definition of $`\theta _n`$ and the fact that $`S_n(v_1)=v_n=x_nx_{n1}`$ for all $`n`$, we have that
$`p(t)={\displaystyle _{\mathrm{}}^t}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\theta _n(s)S_n(v_1)ds={\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\left({\displaystyle _{\mathrm{}}^t}\theta _n(s)𝑑s\right)v_n`$
$`=\left({\displaystyle _{\mathrm{}}^{\mathrm{}}}\theta (s)𝑑s\right){\displaystyle \underset{n=1}{\overset{N1}{}}}v_n+\left({\displaystyle _{\mathrm{}}^t}\theta _N(s)𝑑s\right)v_N`$
$`\left({\displaystyle _{\mathrm{}}^{\mathrm{}}}\theta (s)𝑑s\right){\displaystyle \underset{n=1}{\overset{N1}{}}}v_n+\left({\displaystyle _{\mathrm{}}^{\mathrm{}}}\theta (s)𝑑s\right)v_N`$
$`=\left({\displaystyle _{\mathrm{}}^{\mathrm{}}}\theta (s)𝑑s\right)\left(x_{N1}x_0+x_Nx_{N1}\right){\displaystyle \frac{3}{2}}(2+2)=6.`$
This shows that the image of $`p`$ is contained in the ball $`B(0,6)`$ and $`p`$ is bounded. Let us also remark that $`2v_1\frac{x_1^{}(x_1x_0)}{x_1^{}}\frac{1}{4}`$.
Finally, let us check that $`p`$ satisfies the separation condition (5). Let $`0t<r`$ and take $`N`$ so that $`N1\frac{1}{2}<rN\frac{1}{2}`$; then we have
$`p(r)p(t)={\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\left({\displaystyle _t^r}\theta _n(s)𝑑s\right)v_n={\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\left({\displaystyle _t^r}\theta _n(s)𝑑s\right)(x_nx_{n1})=`$
$`\left({\displaystyle _t^r}\theta _1(s)𝑑s\right)x_0+{\displaystyle \underset{k=1}{\overset{N1}{}}}\left({\displaystyle _t^r}\theta _k(s)𝑑s{\displaystyle _t^r}\theta _{k+1}(s)𝑑s\right)x_k+\left({\displaystyle _t^r}\theta _N(s)𝑑s\right)x_N.`$
By observing that $`\mathrm{max}\{1s,2s1\}\frac{1}{3}`$ for all $`s`$ and taking into account the definition of the $`\theta _n`$, it is not difficult to see that
$`(6)`$
$$\mathrm{max}\{_t^r\theta _N(s)𝑑s,_t^r\theta _{N1}(s)𝑑s_t^r\theta _N(s)𝑑s\}\mathrm{min}\{\frac{1}{3}|tr|,a\},$$
where $`a=_{\frac{1}{2}}^0\theta (s)𝑑s>0`$. Then, by applying either $`x_N^{}`$ or $`x_{N1}^{}`$ to the expression for $`p(r)p(t)`$ above, depending on which the maximum in (6) is, and bearing in mind that $`x_n^{}(x_k)=\delta _{n,k}`$ and $`x_n^{}4`$ for all $`n,k`$, we get that
$$\mathrm{max}\{x_N^{}(p(r)p(t)),x_{N1}^{}(p(r)p(t))\}\mathrm{min}\{\frac{1}{3}|tr|,a\},$$
and it follows that $`p(r)p(t)\mathrm{min}\{\frac{1}{12}|tr|,\frac{a}{4}\}`$. This shows that if $`R>0`$ is large enough then
$$p(r)p(t)\frac{1}{R}\mathrm{min}\{1,|tr|\}$$
for all $`t,r0`$.
In order to get paths $`\beta `$ and $`p`$ and a vector $`v`$ with properties (1)–(5) and such that $`p`$ is contained in the unit ball, it is enough to multiply them all by $`\frac{1}{6}`$. ∎
Proof of theorem 2.1.
Consider $`X=H[z]=H\times `$ and $`C_\epsilon =\{x+tzX:x<\epsilon ,t>0\}`$, where $`H=\text{Ker}z^{}`$ for some $`z^{}X^{}`$ with $`z^{}(z)=z^{}=z=1`$, and $`\epsilon >0`$ is to be fixed later. Let $`\beta `$ and $`p`$ be the paths from lemma 2.2. There is no loss of generality if we assume that $`v[z]`$, $`z^{}(v)\frac{1}{K}`$. Let us define $`\pi :C_\epsilon X`$ by
$$\pi (x,t)=\beta (t)(x)+p(t).$$
It is clear that $`\pi `$ is $`C^{\mathrm{}}`$ smooth and has a bounded derivative. We are going to show that $`\pi `$ is a diffeomorphim onto its image, $`T_\epsilon `$, and $`\pi ^1:T_\epsilon C_\epsilon `$ has a bounded derivative as well. To this end let us define the path $`\alpha :[0,\mathrm{})X^{}`$ by
$$\alpha (t)=f_t=z^{}\beta ^1(t).$$
This is a $`C^{\mathrm{}}`$ smooth and Lipschitz path in $`X^{}`$, and $`\alpha (t)=f_t`$ satisfies that $`\text{Ker}f_t=\beta (t)(H)`$. It is clear from this definition and the properties of $`\beta `$ and $`p`$ that
* $`\alpha ^{}(t)K`$, and
* $`\alpha (t)(p^{}(t))=z^{}(v)\frac{1}{K}`$
for all $`t0`$. Now, for a fixed (but arbitrary) $`yT_\epsilon =\pi (C_\epsilon )`$, let us introduce the auxiliary function $`F=F_y:[0,\mathrm{})`$ defined by
$$F(t)=\alpha (t)(yp(t)).$$
We have that
$`F_y^{}(r)=\alpha ^{}(r)(yp(r))\alpha (r)(p^{}(r))`$
$`\alpha ^{}(r)yp(r)\alpha (r)(p^{}(r))`$
$`Kyp(r){\displaystyle \frac{1}{K}}`$
for all $`r0`$. If we choose $`\epsilon >0`$ smaller than $`\frac{1}{6K^5}`$ this implies that $`\pi `$ is a $`C^{\mathrm{}}`$ diffeomorphism onto its image.
Indeed, let us first see that $`\pi `$ is an injection. Assume that $`y=\pi (x,t)=\pi (w,s)`$ for some $`(x,t),(w,s)C_\epsilon `$. Then we have $`yp(t)=\beta (t)(x)`$ and $`yp(s)=\beta (s)(w)`$, so that $`x=\beta ^1(t)(yp(t))`$ and $`w=\beta ^1(s)(yp(s))`$, and, in order to conclude that $`(x,t)=(w,s)`$, it is enough to see that $`t=s`$. Note that $`\beta (t)(x)\beta (s)(w)=p(s)p(t)`$ and therefore, by (5) of lemma 2.2,
$`{\displaystyle \frac{1}{K}}\mathrm{min}\{1,|ts|\}p(s)p(t)=\beta (t)(x)\beta (s)(w)`$
$`\beta (t)(x)+\beta (s)(w)K(x+w)2K\epsilon {\displaystyle \frac{1}{3K^4}},`$
so that $`|ts|2K^2\epsilon \frac{1}{3K^3}`$. Now, since $`p`$ and $`\beta `$ are both $`K`$-Lipschitz, for every $`r[t,s]`$ we have that
$`yp(r)yp(t)+p(t)p(r)=\beta (t)(x)+p(t)p(r)`$
$`Kx+K|tr|K\epsilon +2K^3\epsilon 3K^3\epsilon .`$
By combining this with the above estimation for $`F_y^{}(r)`$ we get
$`(7)`$
$$F_y^{}(r)Kyp(r)\frac{1}{K}3K^4\epsilon \frac{1}{K}\frac{1}{2K}$$
for every $`r[t,s]`$. Now suppose that $`ts`$. Then, since $`x=\beta ^1(t)(yp(t))`$ and $`w=\beta ^1(s)(yp(s))`$ are both in $`H`$ we have that $`0=z^{}(x)=z^{}(w)=F_y(t)=F_y(s)`$, so that, by the classic Rolle’s theorem, there should exist some $`r(t,s)`$ with $`F_y^{}(r)=0`$. But this contradicts (7). Therefore $`t=s`$ and $`\pi `$ is an injection.
If, for a given $`y\pi (C_\epsilon )`$, we denote by $`t(y)`$ the unique $`t=t(y)`$ such that $`y=\pi (\beta ^1(t)(yp(t)),t)`$ then it is clear that the inverse $`\pi ^1:T_\epsilon C_\epsilon `$ is defined by
$`(8)`$
$$\pi ^1(y)=(\beta ^1(t(y))(yp(t(y))),t(y)).$$
For each $`y`$ the number $`t(y)`$ is uniquely determined by the equation
$$G(y,t):=F_y(t)=0,$$
and the argument above shows that
$`(9)`$
$$\frac{G}{t}(y,t)=F_y^{}(t)\frac{1}{2K}$$
for every $`yT_\epsilon `$ and $`t`$ in a neighbourhood of $`t(y)`$. Then, according to the implicit function theorem we get that the function $`yt(y)`$ is $`C^{\mathrm{}}`$ smooth. Furthermore, we have that
$$t^{}(y)=\frac{\frac{G}{y}(y,t(y))}{\frac{G}{t}(y,t(y))}=\frac{z^{}\beta ^1(t(y))}{F_y^{}(t(y))},$$
and therefore, according to the above estimations,
$$t^{}(y)z^{}\beta ^1(t(y))\frac{1}{|F_y^{}(t(y))|}2K^2,$$
which shows that $`yt(y)`$ has a bounded derivative as well. Then it is clear that $`\pi ^1`$ is $`C^{\mathrm{}}`$ and has a bounded derivative (all the functions involved in (8) have been proved to have bounded derivatives). This concludes the proof of theorem 2.1.
We will finish this section with a simple alternative proof of the failure of Rolle’s theorem in the non-Lipschitz case.
###### Remark 2.3.
If we drop the Lipschitz condition from the statement of theorem 1.1, a much simpler proof based on the same idea is available. Let us make a sketch of this proof.
Consider the decomposition $`X=H\times `$ and pick a non-negative $`C^p`$ smooth bump function $`\phi `$ on $`H`$ whose support is contained on the ball $`B_H(0,1/16)`$. First, we construct a $`C^{\mathrm{}}`$ smooth path $`q:[0,\mathrm{})B_H`$, where $`B_H`$ stands for the unit ball of the hyperplane $`H`$, with the property that $`q`$ has no accumulation points at the infinity, that is, $`lim_n\mathrm{}q(t_n)`$ does not exist for any $`(t_n)`$ going to $`\mathrm{}`$. This can easily be done by having $`q`$ lost in the infinitely many dimensions of $`H`$. For instance, take a biorthogonal sequence $`\{x_n,x_n^{}\}H\times H^{}`$ so that $`x_n=1`$ and $`x_n^{}4`$, and consider a $`C^{\mathrm{}}`$ function $`\theta :[0,1]`$ so that $`\text{supp}\theta [1,1]`$, $`\theta (0)=1`$, $`\theta ^{}(t)<0`$ for $`t(0,1)`$, and $`\theta (t1)=1\theta (t)`$ for $`t[0,1]`$. The path $`q`$ may be defined as
$$q(t)=\underset{n=1}{\overset{\mathrm{}}{}}\theta (tn+1)x_n$$
for $`t0`$. Now we reparametrize $`q`$ and define $`p:[0,1)B_H`$ by
$$p(t)=q\left(\frac{t}{1t}\right).$$
Let $`\alpha :[0,1]`$ be a $`C^{\mathrm{}}`$ smooth function so that $`\alpha (t)=0`$ for all $`t0`$, and $`\alpha ^{}(t)>0`$ for all $`t>0`$. Then the function $`g:X=H\times `$ defined by
$$g(x,t)=\{\begin{array}{cc}& \phi (xp(t))\alpha (t)\text{if}t[0,1);\hfill \\ & 0\text{otherwise}\hfill \end{array}$$
is a $`C^p`$ smooth bump function which does not satisfy Rolle’s theorem. Indeed, it is easy to see that
$$g^{}(x,t)(p^{}(t),1)=\phi (xp(t))\alpha ^{}(t)>0,$$
and in particular $`g^{}(x,t)0`$, for all $`(x,t)`$ in the interior of the support of $`g`$.
## 3. Killing singularities. The failure of Brouwer’s fixed point theorem in infinite dimensions.
Do not be afraid, this section does not contain any totalitarian propaganda. Here we will present two applications of theorem 1.5, both of which have in common the following principle: if you have a mapping with a single singular point or an isolated set of singularities that bother you, you can just kill them by composing your map with some deleting diffeomorphisms. In this way you obtain a new map which is as close as you want to the old one but does not have the adverse properties created by the singular points you eliminate.
For instance, if you want a smooth bump function $`g`$ which does not satisfy Rolle’s theorem and whose support is a smooth starlike body $`A`$, by composing the Minkowski functional of this body with a real bump function you get a function $`h`$ whose support is $`A`$ and whose derivative vanishes only at the origin and outside $`A`$; then, by composing $`h`$ with a diffeomorphism $`f`$ which extracts a small set containing the origin and which restricts to the identity outside $`A`$, you get a map $`g`$ with the required properties.
On the other hand, suppose you want a smooth retraction $`r`$ from a bounded starlike body $`A`$ of a Banach space $`X`$ onto its boundary $`A`$. This is impossible if $`X`$ is finite-dimensional, but otherwise you can use the following trick: it is trivial that there is a smooth retraction $`h`$ from $`A\{0\}`$ onto $`A`$; then take a diffeomorphism $`f`$ which removes from $`X`$ a small subset containing the origin and restricts to the identity outside $`A`$. The composition $`r=hf`$ gives the required retraction.
Let us formalize these ideas and comment on the results that they provide.
The support of the bumps that violate Rolle’s theorem.
The bump function constructed in the proof of theorem 1.1 has a weird support, namely a twisted tube. Some readers might judge this fact rather unpleasant and wonder whether it is possible to construct a bump function which does not satisfy Rolle’s theorem and whose support is a nicer set, such as a ball or a starlike body. To comfort those readers let us first recall that in infinite dimensions there is no topological difference between a tube (whether it is twisted or not) and a ball or a starlike body, as theorem 2.4 in shows. Furthermore, as we said above, theorem 1.5 allows us to show that for a given $`C^p`$ smooth bounded starlike body $`A`$ in an infinite-dimensional Banach space $`X`$, it is always possible to construct a $`C^p`$ smooth bump function on $`X`$ which does not satisfy Rolle’s theorem and whose support is precisely the body $`A`$.
###### Theorem 3.1.
Let $`X`$ be an infinite-dimensional Banach space with a $`C^p`$ smooth bounded starlike body $`A`$. Then there exists a $`C^p`$ smooth bump function $`g`$ on $`X`$ whose support is precisely the body $`A`$, and with the property that $`g^{}(x)0`$ for all $`x`$ in the interior of $`A`$ (that is, $`g`$ does not satisfy Rolle’s theorem).
###### Proof.
Let $`q_A`$ be the Minkowski functional of $`A`$. We may assume that $`B_XA`$. By theorem 1.5 there is a closed subset $`D`$ of $`A`$ and a $`C^p`$ diffeomorphism $`f:XXD`$ which is the identity outside $`A`$. It can be assumed that the origin belongs to $`D`$. Then the function $`h:X`$ defined by
$$h(x)=q_A(f(x))$$
is $`C^p`$ smooth on $`X`$, restricts to the gauge $`q_A`$ outside $`A`$, and has the remarkable property that $`h^{}(x)0`$ for all $`xX`$ (indeed, $`h^{}(x)=q_A^{}(f(x))f^{}(x)`$ is non-zero everywhere because $`q_A^{}(y)0`$ whenever $`y0`$, $`0f(X)`$, and $`f^{}(x)`$ is a linear isomorphism at each point $`x`$).
Now, take a $`C^{\mathrm{}}`$ real function $`\theta :[0,1]`$ such that $`\theta (t)>0`$ for $`t(1,1)`$, $`\theta =0`$ outside $`[1,1]`$, $`\theta (t)=\theta (t)`$, $`\theta (0)=1`$, and $`\theta ^{}(t)<0`$ for all $`t(0,1)`$. Then, if we define $`g:X`$ by
$$g(x)=\theta (h(x)),$$
it is immediately checked that $`g`$ is a $`C^p`$ smooth bump on $`X`$ which does not satisfy Rolle’s theorem and whose support is precisely the body $`A`$. ∎
The failure of Brouwer’s fixed point theorem in infinite dimensions.
The celebrated Brouwer’s fixed point theorem tells us that every continuous self-map of the unit ball of a finite-dimensional normed space admits a fixed point. This is the same as saying that there is no continuous retraction from the unit ball onto the unit sphere, or that the unit sphere is not contractible (the identity map on the sphere is not homotopic to a constant map). The Scottish book (see ) contains the following question, asked in 1935 by S. Ulam. Can one transform continuously the solid sphere of a Hilbert space into its boundary such the transformation should be the identity on the boundary of the ball? In other words, is the unit sphere of a Hilbert space a retract of its unit ball? The first answer to this question is commonly attributed to S. Kakutani , who solved the problem by exhibiting several examples of continuous self-mappings of the unit ball of the Hilbert space without fixed points. Thus, none of the above forms of Brouwer’s fixed point theorem remains valid in infinite dimensions.
A very nice solution to the retraction problem, and one which has the advantage of holding in arbitrary infinite-dimensional Banach spaces, was given by the pioneering results of Klee’s on topological negligibility of points and compacta : for every infinite-dimensional Banach space $`X`$ there always exists a homeomorphism $`h:XX\{0\}`$ so that $`h`$ restricts to the identity outside the unit ball $`B_X`$. The required retraction of $`B_X`$ onto the unit sphere $`S_X`$ is then given by $`R(x)=h(x)/h(x)`$ for $`xB_X`$. By taking into account the subsequent progress on topological negligibility of subsets made by C. Bessaga, T. Dobrowolski and the first-named author among others (see ), this mapping $`h`$ may even be assumed $`C^p`$ smooth provided that the sphere $`S_X`$ is $`C^p`$ smooth. Thus, there are very regular retractions that provide an answer to Ulam’s question.
In B. Nowak showed that for several infinite-dimensional Banach spaces Brouwer’s theorem fails even for Lipschitz mappings (that is, under the strongest uniform-continuity condition), and in Y. Benyamini and Y. Sternfeld generalized Nowak’s result for all infinite-dimensional normed spaces, establishing that for every infinite-dimensional space $`(X,)`$ there exists a Lipschitz retraction from the unit ball $`B_X`$ onto the sphere $`S_X`$, and that $`S_X`$ is Lipschitz contractible.
More recently, M. Cepedello and the first-named author showed that these results hold for the smooth category as well (see ). In fact they proved that for every infinite-dimensional Banach space with a $`C^p`$ (Lipschitz) bounded starlike body $`A`$ (where $`p=0,1,2,\mathrm{},\mathrm{}`$), there is a $`C^p`$ Lipschitz retraction of $`A`$ onto its boundary $`A`$, the boundary $`A`$ is $`C^p`$ (Lipschitz) contractible, and there is a $`C^p`$ smooth (Lipschitz) mapping $`f:AA`$ such that $`f`$ has no (approximate) fixed points.
The proof of these results in the general case is somewhat involved, but if we drop the Lipschitz condition then the fact that Brouwer’s theorem is false in infinite dimensions even for smooth self-mappings of balls or starlike bodies is a trivial consequence of theorem 1.5.
###### Corollary 3.2 (Azagra–Cepedello).
Let $`X`$ be an infinite-dimensional Banach space and let $`A`$ be a $`C^p`$ smooth bounded starlike body. Then:
* The boundary $`A`$ is $`C^p`$ contractible.
* There is a $`C^p`$ smooth retraction from $`A`$ onto $`A`$.
* There exists a $`C^p`$ smooth mapping $`\phi :AA`$ without fixed points.
###### Proof.
Let $`f:XXD`$ be the diffeomorphism from theorem 1.5. We may assume that the origin belongs to the deleted set $`D`$ and that $`B_XA`$, so that $`f`$ restricts to the identity outside $`A`$. Then the formula
$$R(x)=\frac{f(x)}{q_A(f(x))},$$
where $`q_A`$ is the Minkowski functional of $`A`$, defines a $`C^p`$ smooth retraction from $`A`$ onto the boundary $`A`$. This proves (2).
Once we have such a retraction it is easy to prove parts (1) and (3): the formula $`\phi (x)=R(x)`$ defines a $`C^p`$ smooth self-mapping of $`A`$ without fixed points. On the other hand, if we pick a non-decreasing $`C^{\mathrm{}}`$ function $`\zeta :`$ so that $`\zeta (t)=0`$ for $`t\frac{1}{4}`$ and $`\zeta (t)=1`$ for $`t\frac{3}{4}`$, then the formula
$$H(t,x)=R((1\zeta (t))x),$$
for $`t[0,1]`$, $`xA`$, defines a $`C^p`$ homotopy joining the identity to a constant on $`A`$, that is, $`H`$ contracts the pseudosphere $`A`$ to a point. ∎
## 4. How small can the set of gradients of a bump be?
If $`b`$ is a smooth bump function on a Banach space $`X`$ it is natural to ask how large or how small the cone generated by the set of gradients $`b^{}(X)`$ can be. In general, as a consequence of Ekeland’s variational principle, one has that the cone $`𝒞(b)=\{\lambda b^{}(x):\lambda 0,xX\}`$ is norm-dense in the dual space $`X^{}`$ (see , pag. 58, proposition 5.2).
In a study was initiated on the topological properties of the set of derivatives of smooth functions. Among other results it was proved that an infinite-dimensional separable Banach space has a $`C^1`$ smooth bump function (resp. is Asplund) if and only if there exists another $`C^1`$ smooth bump function $`b`$ on $`X`$ with the property that $`b^{}(X)=X^{}`$. This answers the question as to how large can the cone $`𝒞(b)`$ be. But what is the smallest possible size of $`𝒞(b)`$?
To begin with, by using theorem 1.1, one can easily construct smooth bump functions whose sets of gradients lack not only the point zero, but any pre-set finite-dimensional linear subspace of the dual.
###### Corollary 4.1.
Let $`X`$ be an infinite-dimensional Banach space and $`W`$ a finite-dimensional subspace of $`X^{}`$. The following statements are equivalent.
* $`X`$ has a $`C^p`$ smooth (Lipschitz) bump function.
* $`X`$ has a $`C^p`$ smooth (Lipschitz) bump function $`f`$ satisfying that $`\{\lambda f^{}(x):xX,\lambda \}W=\{0\}`$. Moreover, $`\{f^{}(x):f(x)0\}W=\mathrm{}`$.
###### Proof.
We can write $`X=YZ`$, where $`Y=_{w^{}W}\text{Ker}w^{}`$ and $`\text{dim}Z=\text{dim}W`$ is finite. Pick a $`C^p`$ smooth (Lipschitz) bump function $`\phi :Y`$ such that $`\phi `$ does not satisfy Rolle’s theorem, and let $`\theta `$ be a $`C^{\mathrm{}}`$ smooth Lipschitz bump function on $`Z`$ so that $`\theta ^{}(z)=0`$ whenever $`\theta (z)=0`$. Then the function $`f:X=YZ`$ defined by $`f(y,z)=\phi (y)\theta (z)`$ is a $`C^p`$ smooth (Lipschitz) bump which satisfies $`\{f^{}(x):f(x)0\}W=\mathrm{}`$. Indeed, we have
$$f^{}(y,z)=(\theta (z)\phi ^{}(y),\phi (y)\theta ^{}(z))X^{}=Y^{}Z^{}=Y^{}W$$
and, since $`\theta (z)\phi ^{}(y)0`$ whenever $`\phi (y)\theta ^{}(z)0`$, it is clear that $`f^{}(y,z)W`$ unless $`f^{}(y,z)=0`$. ∎
We also have the following
###### Corollary 4.2.
Let $`X`$ be an infinite-dimensional Banach space such that $`X=X_1X_2`$, where $`X_1`$ and $`X_2`$ are both infinite-dimensional. The following statements are equivalent.
* $`X`$ has a $`C^p`$ smooth (Lipschitz) bump function.
* $`X`$ has a $`C^p`$ smooth (Lipschitz) bump function $`f`$ satisfying that $`\{\lambda f^{}(x):xX,\lambda \}(X_1^{}X_2^{})=\{0\}`$. Moreover, $`\{f^{}(x):f(x)0\}(X_1^{}X_2^{})=\mathrm{}`$.
###### Proof.
The proof is similar to that of the preceding corollary. Pick $`\phi _1`$ and $`\phi _2`$ smooth (Lipschitz) bump functions on $`X_1`$ and $`X_2`$, respectively, so that $`\phi _1`$ and $`\phi _2`$ do not satisfy Rolle’s theorem. Then the function $`f:X=X_1X_2`$ defined by $`f(x_1,x_2)=\phi _1(x_1)\phi _2(x_2)`$ is a smooth (Lipschitz) bump which satisfies $`\{f^{}(x):f(x)0\}(X_1^{}X_2^{})=\mathrm{}`$. ∎
If we restrict the scope of our search to classic Banach spaces, much stronger results are available. On the one hand, if $`X=c_0`$ the size of $`𝒞(b)`$ can be really small. Indeed, as a consequence of P. Hájek’s work on smooth functions on $`c_0`$ we know that if $`b`$ is $`C^1`$ smooth with a locally uniformly continuous derivative (note that there are bump functions with this property in $`c_0`$), then $`b^{}(X)`$ is contained in a countable union of compact sets in $`X^{}`$ (and in particular has empty interior). On the other hand, if $`X`$ is non-reflexive and has a separable dual, there are bumps $`b`$ on $`X`$ so that $`𝒞(b)`$ has empty interior, as it was shown in .
In the reflexive case, however, the problem is far from being settled. To begin with, the cone $`𝒞(b)`$ cannot be very small, since it is going to be a residual subset of the dual $`X^{}`$. Indeed, as a straightforward consequence of Stegall’s variational principle, for every Banach space $`X`$ having the Radon-Nikodym Property (RNP) it is easy to see that $`𝒞(b)`$ is a residual set in $`X^{}`$. Therefore, for infinite-dimensional Banach spaces $`X`$ enjoying RNP one can hardly expect a better answer to the above question than the following one: there are smooth bumps $`b`$ on $`X`$ such that the cone $`𝒞(b)`$ has empty interior in $`X^{}`$. In the case of the Hilbert space $`X=\mathrm{}_2`$ we next show that this indeed happens.
###### Theorem 4.3.
In the Hilbert space $`\mathrm{}_2`$ the following holds:
* The usual norm $`||||_2`$ can be uniformly approximated by $`C^1`$ Lipschitz functions $`\psi `$ (with Lipschitz derivative) so that the cones $`𝒞(\psi )`$ generated by the sets of derivatives of $`\psi `$ have empty interior, and $`\psi ^{}(x)0`$ for all $`x\mathrm{}_2`$.
* There is a $`C^1`$ Lipschitz bump function $`b`$ (with Lipschitz derivative) on $`\mathrm{}_2`$ satisfying that the cone $`𝒞(b)`$ generated by its set of derivatives $`b^{}(\mathrm{}_2)`$ has empty interior, and $`b^{}(x)0`$ for every $`x`$ in the interior of its support.
###### Proof.
To save notation, let us just write $`||||`$ when referring to the usual norm in $`\mathrm{}_2`$. We will make use of the following restatement of a striking result due to S. A. Shkarin (see ).
###### Theorem 4.4 (Shkarin).
There is a $`C^{\mathrm{}}`$ diffeomorphism $`\phi `$ from $`\mathrm{}_2`$ onto $`\mathrm{}_2\{0\}`$ such that all the derivatives $`\phi ^{(n)}`$ are uniformly continuous on $`\mathrm{}_2`$, and $`\phi (x)=x`$ for $`x1`$.
Let us consider, for $`0<\epsilon <1`$, the difeomorphism $`\phi _\epsilon :\mathrm{}_2\mathrm{}_2\{0\}`$, $`\phi _\epsilon (x)=\epsilon \phi (x/\epsilon )`$, and the function $`U=U_\epsilon :\mathrm{}_2`$ defined by $`U(x)=\epsilon ^2+\phi _\epsilon (x)^2`$. Then $`U`$ satisfies the following properties:
* $`U`$ is $`C^{\mathrm{}}`$ smooth.
* $`x^2U(x)2\epsilon ^2+x^2`$ and $`\epsilon ^2U(x)`$, for every $`x\mathrm{}_2`$.
* $`U(x)=\epsilon ^2+x^2`$, for every $`x\mathrm{}_2`$, $`x\epsilon `$.
* $`U^{}(x)0`$ for every $`x\mathrm{}_2`$.
* $`U`$ is Lipschitz in bounded sets and $`U^{}`$ is Lipschitz.
Now, we define the functions $`U_n:\mathrm{}_2`$ by $`U_n(x)=\frac{1}{2^{2n}}U(2^nx)`$, whenever $`x\mathrm{}^2`$. We will identify $`\mathrm{}_2`$ with the infinite sum $`_2\mathrm{}_2\mathrm{}_2_2\mathrm{}_2_2\mathrm{}_2\mathrm{}`$, where an element $`x=(x_n)`$ belongs to $`_2\mathrm{}_2`$ if and only if every $`x_n`$ is in $`\mathrm{}_2`$ and $`_nx_n^2<\mathrm{}`$, being $`x^2=_nx_n^2`$. Then, we define the function $`f:_2\mathrm{}_2`$ by
$$f(x)=\underset{n}{}U_n(x_n),\text{ where }x=(x_n)_n.$$
First, note that $`f`$ is well-defined, since condition (ii) implies that, whenever $`x=(x_n)_2\mathrm{}_2`$,
$`0<f(x)={\displaystyle \underset{n}{}}{\displaystyle \frac{1}{2^{2n}}}U(2^nx_n){\displaystyle \underset{n}{}}{\displaystyle \frac{1}{2^{2n}}}(2\epsilon ^2+2^nx_n^2)`$
$`={\displaystyle \underset{n}{}}({\displaystyle \frac{2\epsilon ^2}{2^{2n}}}+||x_n||^2)<\mathrm{}.(4.1)`$
On the one hand, note that, if $`U^{}`$ has Lipschitz constant less than or equal to $`M`$ then $`U_n^{}`$ is also Lipschitz with constant less than or equal to $`M`$, since for $`x`$ and $`y`$ in $`\mathrm{}_2`$ we have
$$U_n^{}(x)U_n^{}(y)=\frac{1}{2^n}U^{}(2^nx)U^{}(2^ny)\frac{1}{2^n}M2^nxy.$$
This implies that, if $`x=(x_n)_2\mathrm{}_2`$, the functionals $`U_n^{}(x_n)\mathrm{}_2`$ satisfy that $`(U_n^{}(x_n))_n_2\mathrm{}_2`$. Indeed, we have $`U_n^{}(x_n)U_n^{}(0)Mx_n`$, and therefore $`_nU_n^{}(x_n)U_n^{}(0)^2<\mathrm{}`$. Also, $`(U_n^{}(0))=(\frac{1}{2^n}U^{}(0))(_2\mathrm{}_2)^{}_2\mathrm{}_2`$, and then we get that $`T(x)=(U_n^{}(x_n))`$ also belongs to $`_2\mathrm{}_2`$.
Let us now prove that $`f`$ is $`C^1`$ smooth. For every $`x=(x_n)`$ and $`h=(h_n)`$ in $`_2\mathrm{}_2`$, we can estimate
$`|f(x+h)f(x)T(x)(h)|{\displaystyle \underset{n}{}}|U_n(x_n+h_n)U_n(x_n)U_n^{}(x_n)(h_n)|`$
$`{\displaystyle \underset{n}{}}|U_n^{}(x_n+t_nh_n)(h_n)U_n^{}(x_n)(h_n)|\text{ (for some }0t_n1)`$
$`M{\displaystyle \underset{n}{}}h_n^2=Mh^2.`$
Therefore $`f`$ is Fréchet differentiable and $`f^{}(x)=(\frac{1}{2^n}U^{}(2^nx_n))`$. Moreover, $`f^{}`$ is Lipschitz since $`f^{}(x)f^{}(y)^2=_nU_n^{}(x_n)U_n^{}(y_n)^2M^2_nx_ny_n^2=M^2xy^2`$. This implies, in particular, that $`f`$ is Lipschitz on bounded sets.
Let us check that $`f=f_\epsilon `$ uniformly approximates $`||||^2`$. Indeed, from condition (ii) on $`U`$ and (4.1), we have that, for every $`x=(x_n)_2\mathrm{}_2`$,
$`(4.2)`$
$$\mathrm{max}\{\frac{1}{3}\epsilon ^2,x^2\}f(x)\frac{2}{3}\epsilon ^2+x^2,$$
and then,
$`(4.3)`$
$$0f(x)x^2\frac{2}{3}\epsilon ^2.$$
In order to obtain functions which approximate the norm uniformly in $`\mathrm{}_2`$ let us consider $`\psi =\psi _\epsilon =\sqrt{f_\epsilon }`$. According to inequalities (4.2) and (4.3) we have that
$$0\psi x\frac{2\epsilon ^2}{3(\psi +x)}\frac{2}{\sqrt{3}}\epsilon $$
for any $`x_2\mathrm{}_2`$.
Let us check that $`\psi ^{}`$ is bounded. By equation (4.2) we have, for any $`x_2\mathrm{}_2`$,
$$\psi ^{}(x)=\frac{f^{}(x)}{2\psi (x)}\frac{f^{}(x)f^{}(0)}{2\psi (x)}+\frac{f^{}(0)}{2\psi (x)}\frac{M}{2}+\frac{\sqrt{3}}{2\epsilon }f^{}(0).$$
Consequently, $`\psi `$ is Lipschitz with Lipschitz constant, say N. In a similar way, we obtain that $`\psi ^{}`$ is Lipschitz, since for any $`x`$, $`y`$ in $`_2\mathrm{}_2`$,
$`\psi ^{}(x)\psi ^{}(y)={\displaystyle \frac{f^{}(x)f^{}(y)}{2\psi (x)}}+{\displaystyle \frac{f^{}(y)}{2}}\left({\displaystyle \frac{1}{\psi (x)}}{\displaystyle \frac{1}{\psi (y)}}\right)`$
$`{\displaystyle \frac{1}{2}}{\displaystyle \frac{f^{}(x)f^{}(y)}{\psi (x)}}+{\displaystyle \frac{\psi (y)\psi (x)}{\psi (x)}}{\displaystyle \frac{f^{}(y)}{2\psi (y)}}`$
$`{\displaystyle \frac{\sqrt{3}M}{2\epsilon }}xy+{\displaystyle \frac{\sqrt{3}N^2}{\epsilon }}xy.`$
Finally, note that the set $`\{\lambda f^{}(x)=\lambda (U_n^{}(x_n)):x=(x_n)_2\mathrm{}_2,\lambda >0\}`$ is contained in $`\{z=(z_n)_2\mathrm{}_2:z_n0\text{ for every }n\}`$, which has empty interior in $`_2\mathrm{}_2`$. This concludes the proof of (1).
In order to prove (2), we consider a $`C^{\mathrm{}}`$ function $`\theta :^+`$, $`\theta ^{}(t)<0`$ for $`t(0,1)`$, and supp$`\theta =(0,1]`$. Then, we can define a required bump function as the composition $`b(x)=\theta (f(x))`$. Indeed, on the one hand, $`f(0)\frac{2}{3}\epsilon ^2<1`$ and therefore $`b(0)>0`$. On the other hand, $`f(x)x^21`$, whenever $`x1`$, and hence $`b(x)=0`$ for $`x1`$. The bump function $`b`$ is clearly Lipschitz with Lipschitz derivative since $`\theta `$, $`\theta ^{}`$ and $`f^{}`$ are Lipschitz and $`f`$ is Lipschitz on bounded sets. ∎
It is clear that the proof of the preceding theorem could only be adapted for superreflexive Banach spaces $`X`$ which admit a decomposition as an infinite sum of subspaces isomorphic to $`X`$, and therefore the problem whether in a separable reflexive space there can be bumps whose sets of gradients have empty interior remains open in the general case, though the information that has already been gathered seems to point to a final positive solution.
###### Open Problem 4.5.
Let $`X`$ be a (separable reflexive) infinite-dimensional Banach space which admits a $`C^p`$ smooth (Lipschitz) bump function. Is there another $`C^p`$ smooth (Lipschitz) bump function $`f`$ on $`X`$ so that the cone $`𝒞(f)=\{\lambda f^{}(x):\lambda 0,xX\}`$ has empty interior in $`X^{}\mathrm{?}`$
The proof of theorem 4.3 naturally raises the following question.
###### Open Problem 4.6.
Can every equivalent norm on $`\mathrm{}_2`$ be uniformly approximated by $`C^1`$ Lipschitz functions satisfying that the cone generated by their derivatives has empty interior? Furthermore, is it possible to approximate squared equivalent norms in $`\mathrm{}_2`$ uniformly on bounded sets by real-analytic functions (resp. polynomials) $`\psi `$ such that the cones generated by the sets $`\psi ^{}(\mathrm{}_2)`$ have empty interior?
Acknowledgements
This research was carried out during a postdoctoral stay of the authors in the Equipe d’Analyse de l’Université Pierre et Marie Curie, Paris 6. The authors are indebted to the Equipe d’Analyse and very especially to Gilles Godefroy for their kind hospitality and generous advice.
Departamento de Análisis Matemático, Facultad de Ciencias Matemáticas, Universidad Complutense, 28040 MADRID, SPAIN.
Equipe d’Analyse, Université Pierre et Marie Curie–Paris 6. 4, place Jussieu, 75005 PARIS, FRANCE.
E-mail addresses: daniel@sunam1.mat.ucm.es, azagra@ccr.jussieu.fr,
marjim@sunam1.mat.ucm.es, marjim@ccr.jussieu.fr |
warning/0002/astro-ph0002509.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Massive stars play a crucial role in the evolution of galaxies and the whole Universe, because they are the primary sources of radiative ionization and heating of the diffuse medium, they provide most of the nucleosynthesis products to boost the metal content of galaxies and the intergalactic medium, and they constitute a major supply of kinetic energy for galaxies, both through stellar winds during their quiescent phases and, eventually, in the form of fast ejecta from supernova explosions.
Therefore, it is fundamental to reach a proper understanding of the formation processes, the detailed properties and the evolution of massive stars. Despite the fact that their high luminosities make them “easy” targets for detailed observational studies, many aspects and properties of massive star evolution are far from being fully understood. This is because in any stellar generation, massive stars constitute a small fraction of the newly formed stars (say, less than 1% by number), they are “elusive” in that their lifetimes are very short (say, less than 10–20 million years), and often they are heavily obscured by the parent molecular clouds where they were formed, making even their identification rather cumbersome. As a consequence it is not easy to cover all evolutionary phases with direct observations of a statistically significant sample of objects.
One can overcome these difficulties and gain additional insights by considering phenomena that indirectly can provide hints and clues to the problem. In other words, besides studying individual massive stars, one can look at the effects that these stars have on their environments (e.g., HII regions, circumstellar nebulae, SNRs), and infer from there what the stars were doing in special phases of their evolution (e.g., formation, LBV and pre-SN phases, etc.) that would not be accessible in other ways.
Thus, one can use radio observations of supernovae, which probe the circumstellar material ejected by the progenitor stars several thousand years before explosion, to study the very last phases of their evolution. These phases represent a tiny fraction of a massive star lifetime, $``$0.1%, and, therefore, they are extremely difficult to reveal and study with direct observations. Although this is an interesting aspect, we are not going to review it here, but rather we refer the reader to recent papers (Montes et al. 1998, Panagia et al. 1999, Weiler et al. 1999).
Here, we consider and discuss one particular aspect of the evolution of massive stars, namely their formation and early evolution in high metallicity environments. We will show that observations of HII regions in external galaxies show that the ionization of He is much lower than that of H when the O/H ratio in the gas is appreciably higher than solar. This implies that at high metallicities either very massive stars ($`M>25M_{}`$) do not form, or they never reach their expected ZAMS location.
## 2 Ionized Helium in the Milky Way
There is clear observational evidence in the H II regions of our Galaxy that the fractional abundance of ionized helium n(He<sup>+</sup>)/n(H<sup>+</sup>) is not a monotonic function of the galactocentric radius. Moving outwards from the Galactic Center, the ionized He abundance is found to increase in the inner Galaxy, then it attains a maximum near the solar circle, and finally drops in the outer Galaxy (e.g., Mezger & Wink 1983 and references therein). The negative gradient in the outer galaxy reflects a genuine decrease in the He abundance in the outward direction (e.g., Panagia 1980; Güsten & Mezger 1982). The positive gradient in the inner Galaxy instead is an effect of the radial metallicity gradient which produces a systematic variation of the spectrum of the ionizing radiation (Panagia 1980).
The fractional ionization of helium is extremely sensitive to the most energetic part of the radiation field powering an H II complex. Hence, it can provide valuable information on the presence and the abundance of the most massive ($`m`$ 20 M ) stars, which are responsible for most of the radiation with energy in excess of 24.6 eV. Therefore, it is a powerful tool to study how the details of the star formation process vary in different physical environments. There are several mechanisms through which a higher metallicity lowers the He ionization in an H II region:
$``$ The relative number of He-ionizing photons in the stellar spectrum is reduced because of both a stronger line blanketing in the 200–500 Å wavelength range, and a higher continuum opacity.
$``$ The stellar radius becomes larger and the effective temperature decreases for a star of given mass, because of the increased continuum opacity in the sub-atmospheric layers of the star.
$``$ The upper cut-off of the Initial Mass Function (IMF, $`m_U`$, may be shifted to lower masses (e.g., Kahn 1974; Shields & Tinsley 1976).
$``$ A higher metallicity may induce a steeper IMF (i.e. a larger value of the slope $`\alpha `$ of the IMF $`N(m)m^\alpha `$) at least for $`m>`$ 10 M , where the bulk of the ionizing radiation is produced (e.g., Terlevich & Melnick 1983).
In the first two cases, metallicity acts “directly” on the radiation field of the ionizing star cluster, by modifying the stellar spectra without affecting the star formation processes. In the third and fourth case instead, metallicity acts “indirectly” and the changes in the radiation field result from changes in the properties of the IMF.
Panagia (1980) demonstrated that the combined effects of at least the first three processes are needed to explain the He ionization in the Milky Way. Moreover, these processes appear to account for the observed gradient of the effective temperature of the ionizing radiation inferred from the fitting of theoretical models to observations of low-metallicity objects (Talent 1980; Campbell 1988).
## 3 Ionized Helium in External Galaxies
Considering external galaxies, several authors (e.g., Pagel 1986, Viallefond 1988, Robledo-Rella & Firmani 1991) have suggested that a systematic change of the IMF with metallicity is required by observations. Others (e.g., Fierro, Torres - Peimbert & Peimbert 1986, McGaugh 1991) have come to the opposite conclusion, and the controversy is still open. A thorough assessment of this subject is now possible and necessary.
We have considered a large sample of extragalactic H II regions which provides an extensive coverage of a very wide metallicity range (almost a factor of 100), and includes galaxies with a variety of morphological types and luminosities. Such a sample is in many respects much more homogeneous than any sample of galactic H II regions. All the H II regions observed are large (diameter D $`>`$ 50 pc), tenuous (n<sub>e</sub> $`<`$ 500 cm<sup>-3</sup> as derived from the \[S II\] I(6717)/I(6731) ratio) and must be ionized by large OB associations.
Here, we limit our analysis to data published as of February 1992. (A more complete investigation, including the discussion of data published as of December 1999, is in progress and will be completed soon; Lenzuni and Panagia 2000, in preparation). Thus, our sample currently includes 287 H II regions in 46 spiral and irregular galaxies with positive detections of the \[O II\] lines at 3726 and 3729 Å (usually unresolved), of the \[O III\] lines at 4959 and 5007 Å , and of at least one of the He I lines. Additional observations were also collected for 87 “Blue Compact Galaxies” (BCG’s). None of these objects is resolved into individual H II regions, the observations being relative to the entire galaxy or, possibly, to its central, brightest parts. These galaxies appear to be undergoing a stage characterized by a collective mode of star-formation. Their spectra are heavily dominated by H II region-like emission, hence they can be treated for our purposes as giant, isolated, extragalactic H II regions.
## 4 Analysis and Discussion
Ionized helium abundances are shown in Figure 1 as a function of oxygen abundances, for all of the H II regions in spiral and irregular galaxies and the Blue Compact Galaxies of our sample. The long baseline in metallicity offers a unique opportunity to constrain both the abundance of primordial helium Y<sub>P</sub> and to the $`\mathrm{\Delta }`$Y/$`\mathrm{\Delta }`$Z gradient, thus fully determining the “helium enrichment curve” (HEC) which relates the total abundance of helium to the abundance of oxygen. Considering that evolutionary effects always decrease the He<sup>+</sup>/H<sup>+</sup> ratio because the aging of a stellar cluster results in a softening of the ionizing radiation field, and that young clusters are observationally favoured because they are intrinsically brighter than older clusters, the HEC can be derived by determining the upper envelope of the distribution shown in Figure 1. Among the class of curves $`Y=Y_0+\mathrm{\Delta }Y/\mathrm{\Delta }Z\times Z`$ the best fit to the upper envelope of the observations is obtained for $`Y_0`$ = 0.243 and ($`\mathrm{\Delta }`$Y/$`\mathrm{\Delta }`$Z ) = 3.2 (see dashed curve in Fig. 1). This relation is consistent with the observational results of Pagel et al. (1992) as well as with Maeder’s (1992) theoretical models.
An inspection to Figure 1 reveals that the observed He<sup>+</sup>/H<sup>+</sup> ratio appears to be almost constant up to solar O abundance (log(O/H)$`{}_{}{}^{}+12`$8.8) and then it declines rather quickly for higher metallicities. This is a clear sign that He is progessively less ionized as the O abundance increases, and implies that the mean radiation temperature of the ionizing stars becomes lower than about 38,000 K around log(O/H)$`8.5`$.
We find that mechanisms through which metallicity acts “directly” on the radiation field of the ionizing star cluster are not enough to explain the observed gradient of the He ionization fraction with O abundance. Most of the effect appears instead to be due to “indirect” mechanisms, i.e. a marked deficiency of hot stars with increasing metal abundances. There are at least three possible scenarios to explain this fact:
1. The IMF slope becomes steeper for higher metallicities.
2. The IMF upper cutoff moves to lower masses for higher metallicities.
3. The most massive stars become progressively unable to provide ionizing radiation, either because at high metallicities the remnant of their pre-MS cocoons remains optically thick over most of a star’s lifetime, or because pulsational instabilities prevent the most massive stars from reaching their expected ZAMS surface conditions.
Our model calculations show that varying only the slope of the IMF, i.e. point (1), does not give a satisfactory fit to the data because the resulting ionization decline would be too shallow. On the other hand, point (2), i.e. a systematic variation of the IMF upper mass cut-off with metal abundance ($`m_U`$ Z, $`\beta >0`$) can reproduce the observed trend of the He ionization, with $`m_U`$ = 48 M and $`\beta `$ = 0.60. Point (3) could also account for the observations provided that the invoked effects are indeed capable to produce the sharp decline of He ionization as observed.
From the observational point of view there are no direct studies to conclusively discriminate between points (2) and (3). Observations of massive stars near the Galactic Center, such as the Pistol star, the Sickle and the Quintuplet clusters (e.g., Figer et al. 1999 and references therein) seem to favor the third possibility, because they are so bright ($`log(L/L_{})>6`$) that they must be quite massive. On the other hand, one may argue that those clusters are so close to the Galactic Center that tidal forces may drastically affect the dynamical processes that lead to the formation of stars and, therefore, they may not be representative of normal situations. The ideal investigation to clarify this issue should include nebular and stellar spectroscopy of a large sample of HII regions in galaxies which display marked effects of incomplete He ionization, such as M51 or M83. Another discriminant between hypotheses (2) and (3) is that if a lowering of the IMF upper cutoff is the explanation (i.e. point (2)), then the frequency of Wolf-Rayet stars relative to early type stars is expected to be abnormally low in the high-Z regions because in this case the reduced ionization is entirely due to the lack of truly massive stars that are expected to end up as WR stars in their final stages of evolution. |
warning/0002/hep-th0002103.html | ar5iv | text | # 1 Introduction
## 1 Introduction
On analyzing the effect of a non-vanishing B-field on the modular invariant partition function of an orbifold, Vafa realised that there is an ambiguity related to phase factors associated to the contributions of the different twisted sectors of the partition function. It is possible to weight the contributions of the $`(g,h)`$-twisted sector, where $`g,h`$ are elements of the orbifold group $`\mathrm{\Gamma }`$, by a phase $`\beta _{g,h}`$ without spoiling modular invariance. These phases are related to elements of $`H^2(\mathrm{\Gamma },U(1))`$ in such a way that if $`\alpha _{g,h}`$ is a 2-cocycle, the phases are of the form: $`\beta _{g,h}=\alpha _{g,h}\alpha _{h,g}^1`$. These phases simply mean a freedom in changing the phase of some operators: the twisted fields in a sector twisted by $`g`$ pick up a phase $`\beta _{g,h}`$ when the element $`h`$ of the orbifold group is acting on them.
The effect of adding discrete torsion was analyzed in detail in some $`\text{/}\text{\_Z}_N\times \text{/}\text{\_Z}_M`$ heterotic orbifolds with standard embedding . Some of these models have zero or negative Euler characteristic. These models are also related to some complete intersection CY manifolds.
Some examples were discused in detail by the authors of . They found that the twisted fields in models with discrete torsion are related to a deformation of the orbifold singularities (in contrast to a blow-up of the singularities in the case without discrete torsion). However, some of the necessary deformations to completely smooth out the orbifold are absent, i.e. some singularities (conifolds) remain. As was commented by Vafa , the discrete torsion is related to some B-flux on a 2-cycle. The identification of this 2-cycle is carried out in . There, a relation was found between the discrete torsion and the torsion part of the homology of the target space (more precisely, the torsion part of the blow-up of the deformed orbifold).
Discrete torsion is implemented in theories containing $`D`$-branes by using projective representations of the orbifold group . These projective representations are classified by $`H^2(\mathrm{\Gamma },U(1))`$ in complete analogy with the discrete torsion in closed string theories found in . Furthermore, the phases in projective representations appear in the amplitudes as a factor $`\alpha _{g,h}\alpha _{h,g}^1`$ where $`\alpha `$ is the cocycle. Recently the relation between the closed string sector discrete torsion and the discrete torsion of the open string sector has been analyzed in disk amplitudes with a twisted closed state and a photon. In order for this amplitude to be invariant under $`\mathrm{\Gamma }`$ the discrete torsion $`\beta `$ and the 2-cocycle must be related as above. $`D`$-brane charges are in correspondence with irreducible projective representations of $`\mathrm{\Gamma }`$, i.e. for each irreducible representation there is a generator of the charge lattice.
The deformation that was previously analyzed in the context of closed string theories can now be seen as deformations of the superpotential, changing the F-flatness conditions. As in the closed string case, some singularities remain after switching on all the possible deformations (the remaining singularities are again conifolds). The relation of models of this type with the $`AdS/CFT`$ correspondence has been analyzed in .
The B-flux responsible for the discrete torsion parameter is related to some flux over a torsion 2-cycle. A similar problem has been analyzed in Type I strings: B-flux on the tori of an orbifold $`T^4/\text{/}\text{\_Z}_N`$ . There, the new features appear at the level of untwisted tadpoles (reduction of the rank) and at the level of the closed twisted sector (some extra tensor multiplets in six dimensions). These models are T-dual to models with non-commuting Wilson Lines . We will find some of these features also in dealing with the discrete torsion in open string theories.
Another way of understanding the discrete torsion has been realized in . As an analogue of the lift of the action of the orbifold group $`\mathrm{\Gamma }`$ to the gauge group (orbifold Wilson lines), one can understand the discrete torsion as an ambiguity in lifting the orbifold group action to another structure (gerbes). 1-gerbes are related to 2-forms, and the action of the orbifold group on it is described in terms of the group cohomology $`H^2(\mathrm{\Gamma },U(1))`$.
The aim of this article is to understand discrete torsion in orientifold models.<sup>1</sup><sup>1</sup>1A different approach has been taken by the authors of . They discuss a compact type I $`𝐙_2\times 𝐙_2`$ orbifold. Their classification of possible models with unbroken supersymmetry agrees with ours. We restrict ourselves to non-compact orientifold constructions. The compact cases will be treated in a future publication . In section 2, we review some of the characteristic features of discrete torsion in closed string and open string theories . In section 3, we analyze the $`D3`$-brane systems at an orbifold singularity in the presence of discrete torsion . The closed string sector, the tadpole consistency conditions and some models are studied in detail. We analyze also the deformation of the $`\text{/}\text{\_Z}_2\times \text{/}\text{\_Z}_2`$ theory with discrete torsion to get the usual $`\text{/}\text{\_Z}_2`$ and conifold theories . In section 4, non-compact orientifolds are constructed with a set of $`D3`$-branes at the orbifold singularity. In general, some $`D7`$-branes are also needed in order to cancel the Klein bottle contribution to the tadpoles. Some consistency conditions are found in different sectors: the closed string sector, the open string sector and the tadpole cancellation conditions. In all the models we consider, these three conditions lead to the same restriction on the discrete torsion parameter $`ϵ`$: only real values are allowed in the orientifold case, i.e. $`ϵ=\pm 1`$. The orientifold involution is related to real projective representations. Four types of orientifold models are found. This classification is based on the possibility of having vector structure or no vector structure for each of the two generators of $`\text{/}\text{\_Z}_N\times \text{/}\text{\_Z}_M`$. The open string spectrum and the tadpole conditions are different for each of the four cases. We determine explicitly the open string spectrum of the $`D3`$-branes. Some of the deformations of the superpotential, that in the orbifold case are allowed, are not present in the orientifold case. In particular, only the deformation to the $`\text{/}\text{\_Z}_2`$ theory survives in the $`\text{/}\text{\_Z}_2\times \text{/}\text{\_Z}_2`$ model.
Finally, some tools used in this paper are explained in the appendices. Many of the results of this article are based on the theory of of complex and real projective representations . In appendix A, we therefore give a summary of their basic properties. Appendix B contains a detailed computation of the tadpole conditions for $`\text{/}\text{\_Z}_N\times \text{/}\text{\_Z}_M`$ orientifolds. The relation between the shift formalism and the usual construction using matrices to determine the spectrum of orientifold models is explained in appendix C. We show how the shift formalism can be generalized to treat orientifolds with discrete torsion.
## 2 Discrete torsion
In the original paper of Vafa discrete torsion appeared as phase factors in the one-loop partition function of closed string orbifold theories. If we denote by $`Z(g,h)`$ the contribution of the $`(g,h)`$-twisted sector<sup>2</sup><sup>2</sup>2By this we mean the contribution from world-sheets that are twisted by $`g`$ in their space direction and by $`h`$ in their time direction., then the total partition function $`Z`$ is given by
$$Z=\frac{1}{|\mathrm{\Gamma }|}\underset{g,h\mathrm{\Gamma }}{}\beta _{g,h}Z(g,h),$$
(2.1)
where $`\mathrm{\Gamma }`$ is the orbifold group. The usual case without discrete torsion corresponds to $`\beta _{g,h}=1g,h`$. The discovery of Vafa was that non-trivial phases are consistent with modular invariance if they satisfy the following conditions:
$$\beta _{g,g}=1,\beta _{g,h}=\beta _{h,g}^1,\beta _{g,hk}=\beta _{g,h}\beta _{g,k}g,h,k\mathrm{\Gamma }.$$
(2.2)
This implies (see appendix A) that these phases are of the form $`\beta _{g,h}=\alpha _{g,h}\alpha _{h,g}^1`$, where $`\alpha `$ is a 2-cocycle of the group $`\mathrm{\Gamma }`$. The possible discrete torsions are therefore classified by the group cohomology $`H^2(\mathrm{\Gamma },U(1))`$. We will only consider Abelian orbifold groups, i.e. $`\mathrm{\Gamma }`$ is (isomorphic to) a product of cyclic groups. Moreover, as the internal space of the string models we want to discuss is complex three-dimensional, all the possible cases can be reduced to $`\mathrm{\Gamma }=\text{/}\text{\_Z}_N`$ or $`\mathrm{\Gamma }=\text{/}\text{\_Z}_N\times \text{/}\text{\_Z}_M`$. It is known (see e.g. ) that $`H^2(\text{/}\text{\_Z}_N,U(1))=1`$ and $`H^2(\text{/}\text{\_Z}_N\times \text{/}\text{\_Z}_M,U(1))=\text{/}\text{\_Z}_{\mathrm{gcd}(N,M)}`$. As a consequence, discrete torsion is only possible in $`\text{/}\text{\_Z}_N\times \text{/}\text{\_Z}_M`$ orbifolds and is generated by one element of $`\text{/}\text{\_Z}_{\mathrm{gcd}(N,M)}`$. More precisely, let $`g_1,g_2`$ be the generators of $`\text{/}\text{\_Z}_N,\text{/}\text{\_Z}_M`$ respectively, $`p=\mathrm{gcd}(N,M)`$ and choose $`\beta _{g_1,g_2}=\omega _p^m`$, with $`\omega _p`$ a $`p`$-th primtive root of unity and $`m=1,\mathrm{},p`$. All the other phases $`\beta _{g,h}`$ are then fixed by (2.2). In the notation of appendix A, eq. (A.9), they read
$$\beta _{(a,b),(a^{},b^{})}=ϵ^{ab^{}ba^{}},\mathrm{with}ϵ=e^{2\pi im/p},m=1,\mathrm{},p.$$
(2.3)
As the partition function encodes the spectrum of the considered orbifold model, it is clear that discrete torsion modifies the spectrum. From a geometrical point of view this change in the spectrum can be understood from the fact that discrete torsion changes the cohomology of the internal space . However, the spectrum of the untwisted sector remains unchanged, as can be seen from $`\beta _{g,e}=1`$ (this follows from the third equation in (2.2) by setting $`h=g`$ and $`k=e`$, where $`e`$ is the neutral element of $`\mathrm{\Gamma }`$).
The generalization of discrete torsion to open strings has been found by Douglas . The matrices $`\gamma _g`$ that represent the action of the elements $`g`$ of the orbifold group $`\mathrm{\Gamma }`$ on the Chan-Paton indices of the open strings form a projective representation: $`\gamma _g\gamma _h=\alpha _{g,h}\gamma _{gh}`$. The $`\alpha _{g,h}`$ are arbitrary non-zero complex numbers. They are called the factor system of the projective representation $`\gamma `$ and they form a 2-cocycle in the sense that they satisfy (A.2). Two matrices $`\gamma _g`$ and $`\widehat{\gamma }_g`$ are considered projectively equivalent if there exists a non-zero complex number $`\rho _g`$ such that $`\widehat{\gamma }_g=\rho _g\gamma _g`$. As shown in appendix A the set of equivalence classes of cocycles $`\alpha `$ is $`H^2(\mathrm{\Gamma },U(1))`$. Thus, the ambiguity due to the projective representations in the open string sector is classified by the same group cohomology as the discrete torsion of in the closed string sector. It is therefore natural to assume that choosing a non-trivial factor system $`\alpha _{g,h}`$ corresponds to discrete torsion in the open string sector. Moreover it was shown in that a factor system $`\alpha _{g,h}`$ in the open string sector of some orbifold model leads to phases $`\beta _{g,h}=\alpha _{g,h}\alpha _{h,g}^1`$ in the closed string partition function of this model.
## 3 Non-compact orbifold construction
Let us consider a set of $`D3`$-branes at an orbifold singularity of the non-compact space $`\mathrm{C}\text{ }^3/\mathrm{\Gamma }`$, where $`\mathrm{\Gamma }=\text{/}\text{\_Z}_N\times \text{/}\text{\_Z}_M`$. The effect of discrete torsion in such models has been studied in . For completeness and to fix our notations we first summarize their results and then expand on them.
### 3.1 Closed string spectrum
The closed string spectrum in $`D=4`$ can be obtained from the cohomology of the internal space $`\mathrm{C}\text{ }^3/\mathrm{\Gamma }`$. Strictly speaking, this spectrum is continuous because $`\mathrm{C}\text{ }^3/\mathrm{\Gamma }`$ is non-compact. However, as noted in footnote 2 of , it is interesting to determine the massless spectrum that would emerge, if the internal space were compactified. This is what we want to do. The analysis is similar to the one performed by Vafa and Witten . The cohomology can be split in untwisted and twisted contributions: $`H^{p,q}=H_{\mathrm{untw}}^{p,q}+_{g\mathrm{\Gamma }\{e\}}H_g^{p,q}`$. The untwisted Hodge number $`h_{\mathrm{untw}}^{p,q}=dim(H_{\mathrm{untw}}^{p,q})`$ is just given by the number of $`\mathrm{\Gamma }`$-invariant $`(p,q)`$-forms on $`\mathrm{C}\text{ }^3`$. This result is independent of the discrete torsion. The twisted contributions are due to the singularities of $`\mathrm{C}\text{ }^3/\mathrm{\Gamma }`$. In the sector twisted by $`g`$, one has to find the $`\mathrm{\Gamma }`$-invariant forms that can be defined on the subspace $`_g`$ of $`\mathrm{C}\text{ }^3`$ that is fixed under the action of $`g`$. If discrete torsion is present, then the forms on $`_g`$ must be invariant under the action of $`h`$ combined with a multiplication by $`\beta _{g,h}`$ $`h\mathrm{\Gamma }`$. In the non-compact case that we are treating here, $`_g`$ is either a point at the origin of $`\mathrm{C}\text{ }^3`$ or a complex plane located at the origin of the transverse $`\mathrm{C}\text{ }^2`$. In the former case there is only a $`(0,0)`$-form. If no discrete torsion is present, it is invariant and contributes one unit to $`h_g^{1,1}`$ or to $`h_g^{2,2}`$ depending on the specific model. In the case with discrete torsion, there is no contribution to the cohomology from this sector. If $`_g`$ is a complex plane, then, in the case without discrete torsion, the $`(0,0)`$-form and the $`(1,1)`$-form are invariant. They contribute one unit to $`h_g^{1,1}`$ and to $`h_g^{2,2}`$. In the case with discrete torsion the $`(0,0)`$-form and the $`(1,1)`$-form are invariant if $`\beta _{g,h}=1h`$. The $`(0,1)`$-form and the $`(1,0)`$-form are invariant if the action of $`h`$ multiplies these forms by a phase which is opposite to $`\beta _{g,h}h`$. In this case, they contribute one unit to $`h_g^{1,2}`$ and to $`h_g^{2,1}`$.
For simplicity let us first take $`\mathrm{\Gamma }=\text{/}\text{\_Z}_2\times \text{/}\text{\_Z}_2`$ and then indicate the generalization to $`\text{/}\text{\_Z}_N\times \text{/}\text{\_Z}_M`$. Using the method explained above, the Hodge diamond of the untwisted cohomology is found to be
$$\begin{array}{cccccccc}& & & & 1& & & \\ & & & 0& & 0& & \\ & & 0& & 3& & 0& \\ & 1& & 3& & 3& & 1\\ & & 0& & 3& & 0& \\ & & & 0& & 0& & \\ & & & & 1& & & \end{array}$$
For all the other $`\text{/}\text{\_Z}_N\times \text{/}\text{\_Z}_M`$ models, one finds nearly the same untwisted cohomology, the only difference being that in those cases $`h_{\mathrm{untw}}^{2,1}=h_{\mathrm{untw}}^{1,2}=0`$.
There are three twisted sectors in the $`\text{/}\text{\_Z}_2\times \text{/}\text{\_Z}_2`$ corresponding to the three non-trivial elements (0,1), (1,0), (1,1), where we used the notation introduced in appendix A (below eq. (A.8)). Each of these sectors gives the same contribution to the twisted cohomology. In the case without discrete torsion:
$$\begin{array}{cccccccc}& & & & 0& & & \\ & & & 0& & 0& & \\ & & 0& & 1& & 0& \\ & 0& & 0& & 0& & 0\\ & & 0& & 1& & 0& \\ & & & 0& & 0& & \\ & & & & 0& & & \end{array}$$
In the case with discrete torsion:
$$\begin{array}{cccccccc}& & & & 0& & & \\ & & & 0& & 0& & \\ & & 0& & 0& & 0& \\ & 0& & 1& & 1& & 0\\ & & 0& & 0& & 0& \\ & & & 0& & 0& & \\ & & & & 0& & & \end{array}$$
For a general $`\text{/}\text{\_Z}_N\times \text{/}\text{\_Z}_M`$ orbifold one finds that, only if minimal discrete torsion $`ϵ=\omega _{\mathrm{gcd}(N,M)}`$ (i.e. $`m=1`$ in (2.3)) is present, is there a non-zero contribution to $`h_{\mathrm{tw}}^{1,2}=_{g\mathrm{\Gamma }\{e\}}h_g^{1,2}`$. Three cases can be distinguished:
i) $`N=Mh_{\mathrm{tw}}^{1,2}=h_{\mathrm{tw}}^{2,1}=3`$
ii) $`N=\mathrm{gcd}(N,M)<Mh_{\mathrm{tw}}^{1,2}=h_{\mathrm{tw}}^{2,1}=2`$
iii) $`N\mathrm{gcd}(N,M)Mh_{\mathrm{tw}}^{1,2}=h_{\mathrm{tw}}^{2,1}=0`$
In order to obtain the closed string spectrum in $`D=4`$, one has to dimensionally reduce the massless spectrum of type IIB supergravity in $`D=10`$: the metric $`g^{(10)}`$, the NSNS 2-form $`B^{(10)}`$, the dilaton $`\varphi ^{(10)}`$, the RR-forms $`C_0^{(10)}`$, $`C_2^{(10)}`$, $`C_4^{(10)}`$. (We only give the bosons, the fermions are related to them by supersymmetry.) This is done by contracting their Lorentz indices with the differential forms of the internal space. The resulting spectrum has $`𝒩=2`$ supersymmetry in $`D=4`$. For a general configuration we get:
* $`g^{(4)}`$, $`(h^{1,1}+2h^{2,1})`$ scalars (from $`g^{(10)}`$)
* a 2-form, $`h^{1,1}`$ scalars (from $`B^{(10)}`$)
* a scalar (from $`\varphi ^{(10)}`$)
* a scalar (from the $`C_0^{(10)}`$)
* a 2-form, $`h^{1,1}`$ scalars (from $`C_2^{(10)}`$)
* $`(h^{2,1}+1)`$ vectors, $`h^{1,1}`$ 2-forms (from $`C_4^{(10)}`$).
These fit into the following $`𝒩=2`$ SUSY representations (see e.g. ):
* a gravity multiplet (consisting of $`g^{(4)}`$, a vector and fermions)
* a double tensor multiplet (consisting of two 2-forms, two scalars and fermions)
* $`h^{1,1}`$ tensor multiplets (consisting of a 2-form, three scalars and fermions)
* $`h^{2,1}`$ vector multiplets (consisting of a vector, two scalars and fermions)
For the $`\text{/}\text{\_Z}_2\times \text{/}\text{\_Z}_2`$ example we find (in $`𝒩=2`$ multiplets):
\- without discrete torsion: gravity, a double tensor, 3 vectors and 6 tensors
\- with discrete torsion: gravity, a double tensor, 6 vectors and 3 tensors
### 3.2 Tadpoles
In the orbifold case, only oriented surfaces appear in computing the tadpoles. There is only one involving branes: the cylinder. As the space is non-compact all the tadpoles with an inverse dependence on the volume of the internal coordinates vanish. These are the untwisted ones and those corresponding to twisted sectors that leave one complex plane fixed. So there is no restriction on the total number of branes and on $`Tr\gamma _g`$ if $`g`$ has a fixed plane. This means that, for example, in the $`\text{/}\text{\_Z}_2\times \text{/}\text{\_Z}_2`$ case there is no restriction from tadpoles (that is the case of ref. ).
In general, one can add $`D7_i`$-branes to the system of $`D3`$-branes at the singularity, where the index $`i=1,2,3`$ refers to the complex plane with Dirichlet conditions. Then the cylinder amplitude can be split into four sectors: $`33`$, $`7_i7_i`$, $`7_i7_j`$ and $`7_i3`$. It is convenient to denote a group element $`g=g_1^ng_2^m`$ of $`\mathrm{\Gamma }=\text{/}\text{\_Z}_N\times \text{/}\text{\_Z}_M`$ by the 2-vector $`\overline{k}=(n,m)`$. Similarly, we form 2-vectors of the components of the twist vectors $`v_i`$ and $`w_i`$ that represent the action of $`g_1`$ and $`g_2`$ on the $`i`$-th complex plane: $`\overline{v}_i=(v_i,w_i)`$. If we define $`s_i=\mathrm{sin}(\pi \overline{k}\overline{v}_i)`$, then the cylinder contribution is of the form (see appendix B):
$$𝒞=\underset{\overline{k}=(1,1)}{\overset{(N,M)}{}}\frac{1}{8s_1s_2s_3}\left[8s_1s_2s_3Tr\gamma _{\overline{k},3}+_{i=1}^32s_iTr\gamma _{\overline{k},7_i}\right]^2.$$
(3.1)
As $`D7`$-branes are not needed for the consistency of the orbifold models, we restrict ourselves to configurations without $`D7`$-branes. In this case the tadpole conditions are summarized by:
$$Tr\gamma _g=0,\text{if }g\text{ has no fixed planes.}$$
(3.2)
As we have seen, discrete torsion appears in two ways: as phase factors in the closed string partition function and as a non-trivial factor system in the projective representation of the orbifold group on the Chan-Paton indices of open strings. The former could only modify the Klein bottle diagram the latter could only affect the Möbius strip. Both diagrams are not present in the orbifold case. Consequently, in this case, the tadpole conditions are not changed by discrete torsion.
Using formula (A.14) for the characters of a projective representation, it is easy to find solutions to the tadpole conditions (3.2). According to the results of appendix A, the matrix $`\gamma _g`$ (with $`g=g_1^ag_2^b`$) of a general projective representation with discrete torsion $`ϵ`$ is of the form
$$\underset{k,l}{}(\sqrt{ϵ})^{ab}(\omega _N^k\gamma _{g_1})^a(\omega _M^l\gamma _{g_2})^b\mathrm{\hspace{0.17em}1}\mathrm{I}_{n_{kl}},$$
(3.3)
where $`\gamma _{g_1}`$ and $`\gamma _{g_2}`$ are given in (A.11), $`k=0,\mathrm{},N/s1`$, $`l=0,\mathrm{},M/s1`$. We recall that $`\gamma _{g_{1/2}}`$ are $`(s\times s)`$-matrices, where $`s`$ is the smallest positive integer, such that $`ϵ^s=1`$. One can readily verify that the regular representation, i.e. $`n_{kl}=sk,l`$, is a solution of (3.2). But there are many more solutions. For example each set of $`n_{kl}`$ that either only depends on $`k`$ or only depends on $`l`$ (i.e. $`n_{kl}=\stackrel{~}{n}_kl`$ or $`n_{kl}=\stackrel{~}{n}_lk`$) is possible.
To make some further statements, we need to specify how the generators $`g_1`$, $`g_2`$ of $`\text{/}\text{\_Z}_N\times \text{/}\text{\_Z}_M`$ act on the three complex coordinates of the internal space. We choose the following action:
$`g_1:`$ $`(z_1,z_2,z_3)(e^{2\pi iv_1}z_1,e^{2\pi iv_2}z_2,e^{2\pi iv_3}z_3)`$
$`g_2:`$ $`(z_1,z_2,z_3)(e^{2\pi iw_1}z_1,e^{2\pi iw_2}z_2,e^{2\pi iw_3}z_3)`$
$`\mathrm{with}`$ $`v={\displaystyle \frac{1}{N}}(1,1,0),w={\displaystyle \frac{1}{M}}(0,1,1).`$ (3.4)
For $`N=M`$, all possible actions of $`g_1`$, $`g_2`$ can be brought to the form (3.2) by permuting the elements of $`\text{/}\text{\_Z}_N\times \text{/}\text{\_Z}_N`$. If $`NM`$, one can choose different (non-equivalent) actions for $`g_1,g_2`$. We restrict ourselves to one alternative:
$`g_1:`$ $`(z_1,z_2,z_3)(e^{2\pi iv_1}z_1,e^{2\pi iv_2}z_2,e^{2\pi iv_3}z_3)`$
$`g_2:`$ $`(z_1,z_2,z_3)(e^{2\pi iw_1^{}}z_1,e^{2\pi iw_2^{}}z_2,e^{2\pi iw_3^{}}z_3)`$
$`\mathrm{with}`$ $`v={\displaystyle \frac{1}{N}}(1,1,0),w^{}={\displaystyle \frac{1}{M}}(1,2,1).`$ (3.5)
To distinguish the two different actions, we will denote the orbifold corresponding to the latter by $`\text{/}\text{\_Z}_N\times \text{/}\text{\_Z}_M^{}`$.
Two interesting cases are (i) $`N=M`$ and (ii) $`M`$ a multiple of $`N`$. In the first case, there is no tadpole condition for $`s=N`$ and $`s=N/2`$ (the only non-vanishing characters correspond to group elements that have fixed planes, see eq. (A.14)). So the first condition arises for $`s=N/3`$, which only for $`N6`$ leads to an orbifold with discrete torsion. The condition that the $`n_{kl}`$ of (3.3) have to satisfy for $`N=M=6`$ and $`s=2`$ is $`_{k,l=0}^2n_{kl}e^{2\pi i(k+2l)/3}=0`$. As shown above, it is easy to find solutions for the $`n_{kl}`$. For higher $`N`$ more conditions of this type have to be satisfied. If $`M`$ is a multiple of $`N`$, then there are non-vanishing characters even in the case of minimal discrete torsion, $`s=N`$: $`Tr\gamma _{(0,cN)}=N_{l=0}^{M/N1}n_{0,l}\omega _M^{lcN}`$, where $`c=1,\mathrm{},M/N1`$. However, the group element $`(0,cN)`$ has a fixed plane if the group action is as in (3.2). Thus, only for the $`\text{/}\text{\_Z}_N\times \text{/}\text{\_Z}_M^{}`$ orbifold do we get a tadpole condition in the case $`s=N`$. A solution to this condition is to take $`n_{0,l}=nl`$. If $`M/N`$ is a prime number, then this solution is unique, else different choices for $`n_{0,l}`$ are possible.
### 3.3 Open string spectrum
Discrete torsion is implemented in the relation between the elements $`\gamma _g`$ of the representation of the group acting on the Chan-Paton matrices :
$$\gamma _g\gamma _h=\alpha _{g,h}\gamma _{gh}.$$
(3.6)
As in the case without discrete torsion, one gets the gauge fields $`\lambda ^{(0)}`$ and matter fields $`\lambda ^{(i)}`$, $`i=1,2,3`$, taking the solutions to the projections:
$$\gamma _g^1\lambda \gamma _g=r(g)\lambda ,$$
(3.7)
where $`r(g)`$ is the matrix that represents the action of $`g`$ on $`\lambda =(\lambda ^{(0)},\lambda ^{(1)},\lambda ^{(2)},\lambda ^{(3)})`$.
Let us first treat the $`\text{/}\text{\_Z}_2\times \text{/}\text{\_Z}_2`$ model of in detail and then show how this is generalized to more complicated models. The only non-trivial phases in (3.6) are
$$\alpha _{g,h}=\alpha _{h,g}=i,$$
(3.8)
where we have denoted the generators of the first and second $`\text{/}\text{\_Z}_2`$ by $`g`$ and $`h`$. This corresponds to a projective representation with discrete torsion $`ϵ=\alpha _{g,h}\alpha _{h,g}^1=1`$. From (3.3), we find that the matrices are of the form:
$`\gamma _e`$ $`=`$ $`1\mathrm{I}_21\mathrm{I}_n,`$ (3.9)
$`\gamma _g`$ $`=`$ $`\sigma _31\mathrm{I}_n,`$
$`\gamma _h`$ $`=`$ $`\sigma _11\mathrm{I}_n,`$
$`\gamma _{gh}`$ $`=`$ $`\sigma _21\mathrm{I}_n,`$
where $`\sigma _i`$ are the Pauli matrices and $`n`$ is an arbitrary parameter (it counts the number of dynamical $`D`$-branes).
The solution to (3.7) is given by
$$\lambda ^{(0)}=1\mathrm{I}_2X,\lambda ^{(i)}=\sigma _iZ_i,$$
(3.10)
where $`X,Z_i`$ are arbitrary $`(n\times n)`$-matrices. This corresponds to gauge group $`U(n)`$ and three adjoint matter fields $`Z_1`$, $`Z_2`$, $`Z_3`$ in $`𝒩=1`$ multiplets. The superpotential can be obtained in the usual way and reads:
$$W=Tr(Z_1Z_2Z_3+Z_2Z_1Z_3).$$
(3.11)
It has to be completed the by the deformations
$$\mathrm{\Delta }W=\underset{i=1}{\overset{3}{}}\zeta _iTrZ_i,$$
(3.12)
where $`\zeta _i`$ are the twisted modes from the closed string spectrum.
For general $`\gamma `$-matrices (3.3), representing the action of $`\text{/}\text{\_Z}_N\times \text{/}\text{\_Z}_M`$ on the open strings, the solution to (3.7) is a gauge group
$$\underset{k=0}{\overset{N/s1}{}}\underset{l=0}{\overset{M/s1}{}}U(n_{kl})$$
(3.13)
with matter in adjoint and bifundamental representations. For some models the tadpole conditions impose restrictions on the numbers $`n_{kl}`$. The spectrum can easiest be obtained from the corresponding quiver diagram, as explained in . A more rigorous way to find the spectrum uses the shift formalism developed in . We explain this method in appendix C. Note that a $`\text{/}\text{\_Z}_N\times \text{/}\text{\_Z}_M`$ orbifold with discrete torsion characterized by the integer<sup>3</sup><sup>3</sup>3Recall that $`s`$ was defined as the smallest positive integer, such that $`ϵ^s=1`$. $`s`$ has the same open string spectrum as the orbifold $`\text{/}\text{\_Z}_{N/s}\times \text{/}\text{\_Z}_{M/s}`$ without discrete torsion. The only difference appears in the superpotential where some of the terms acquire additional phases .
In table 1 we present the explicit solution for some models.
Note that the $`\text{/}\text{\_Z}_2\times \text{/}\text{\_Z}_6^{}`$ example with discrete torsion $`ϵ=1`$ would have non-Abelian gauge anomalies if the gauge group were $`U(n_1)\times U(n_2)\times U(n_3)`$, with $`n_1n_2n_3`$. However, as mentioned in the previous subsection, the requirement of tadpole cancellation implies that $`n_1=n_2=n_3`$.<sup>4</sup><sup>4</sup>4The relation between anomaly freedom and tadpole cancellation has been discussed in . This can be seen as follows. The action of $`\text{/}\text{\_Z}_2\times \text{/}\text{\_Z}_6^{}`$ on the complex coordinates is as in (3.2). There are three twisted sectors (we are not taking into account the inverse because they give the same tadpole cancellation condition) that do not have any fixed plane: $`h`$, $`h^2`$ and $`gh`$, where $`g`$ and $`h`$ are the generators of $`\text{/}\text{\_Z}_2`$ and $`\text{/}\text{\_Z}_6^{}`$. Three tadpoles must be cancelled:
$$Tr\gamma _h=0,Tr\gamma _{gh}=0,Tr\gamma _{h^2}=0.$$
(3.14)
The action of the orbifold group on the Chan-Paton indices can be chosen as
$`\gamma _h`$ $`=`$ $`diag(1\mathrm{I}_{n_1},1\mathrm{I}_{n_1},e^{2\pi i/6}1\mathrm{I}_{n_2},e^{8\pi i/6}1\mathrm{I}_{n_2},e^{4\pi i/6}1\mathrm{I}_{n_3},e^{10\pi i/6}1\mathrm{I}_{n_3}),`$
$`\gamma _g`$ $`=`$ $`\left(\begin{array}{cccccc}0& 1& & & & \\ 1& 0& & & & \\ & & 0& 1& & \\ & & 1& 0& & \\ & & & & 0& 1\\ & & & & 1& 0\end{array}\right).`$ (3.21)
The first two tadpole conditions are inmediatly satisfied but the third one implies:
$$n_1+e^{2\pi i/3}n_2+e^{2\pi i/3}n_3=0.$$
(3.22)
The solution coincides with the condition for non-Abelian gauge anomaly cancellation: $`n_1=n_2=n_3`$.
A similar calculation shows that, in the case of $`\text{/}\text{\_Z}_6\times \text{/}\text{\_Z}_6`$ with discrete torsion $`ϵ=1`$, the gauge group is not of the most general form (3.13). One needs $`n_{1,1}=n_{1,2}=n_{1,3}`$, $`n_{2,1}=n_{2,2}=n_{2,3}`$ and $`n_{3,1}=n_{3,2}=n_{3,3}`$ for tadpole cancellation. The same condition can also be deduced from the requirement of anomaly freedom.
### 3.4 Resolution of the singularities
$`D`$-branes are interesting probes of geometry and topology. World volume theories of $`D`$-branes at a singularity are related to the structure of this singularity. Spacetime can even be understood as a derived concept, emerging from non-trivial moduli spaces of the theory on the $`D`$-brane .
In our case, there is a correspondence between the moduli space of the theory and the transverse space to the $`D3`$-branes. This correspondence allows a map between different theories and different singularities (this has been extensively studied, see e.g. ).
Cases without discrete torsion are related to deformations of the D-flatness conditions. Fayet-Iliopoulos terms are controlled by the twisted sector moduli. Non-trivial values of the Fayet-Iliopoulos parameters can be interpreted as resolutions of the singularities transverse to the branes. Arbitrary values of these terms resolve completely the singularity.
Cases with discrete torsion are related to deformations of the F-flatness conditions . The singularities are not smoothed out by blow-ups but by deformations . However, these singularities cannot be completely smoothed out due to the lack of a sufficient number of twisted states. Some singularities are remaining.
Resolutions and deformations are the two main strategies of desingularization of orbifolds. An analysis of the different desingularizations associated to some orbifold singularities is done in . Only a small number of the possible desingularizations are available in string theory.
In this section we want to find the field theories related to the deformation of the $`\mathrm{C}\text{ }/\text{/}\text{\_Z}_2\times \text{/}\text{\_Z}_2`$ singularity with discrete torsion $`ϵ=1`$. The analysis in the orientifold case will be derived from this one.
From the three adjoint fields one can construct the $`SU(n)`$ invariants. We consider the diagonal $`U(1)`$ as decoupled . These $`SU(n)`$ invariants are of the form:
$`M_{i,j}`$ $`=`$ $`Tr(Z_iZ_j),`$ (3.23)
$`B`$ $`=`$ $`Tr(Z_1[Z_2,Z_3]).`$
They satisfy the following relation:
$$B^2=detM.$$
(3.24)
Using the F-flatness conditions coming from the superpotential (3.11), (3.12) above, one can obtain the following relation between the invariants:
$$det\left(\begin{array}{ccc}M_{11}& \zeta _3& \zeta _2\\ \zeta _3& M_{22}& \zeta _1\\ \zeta _2& \zeta _1& M_{33}\end{array}\right)=B^2.$$
(3.25)
There are three different singularities derived from this case:
* If all the $`\zeta _i`$ are equal to zero the above equation between the invariants is of the form:
$$M_{11}M_{22}M_{33}=B^2.$$
(3.26)
This is a $`\text{/}\text{\_Z}_2\times \text{/}\text{\_Z}_2`$ singularity.
* If only one of the three $`\zeta _i`$, say $`\zeta _1`$, is different from zero, the equation relating the invariants is:
$$M_{11}(M_{22}M_{33}\zeta _1^2)=B^2.$$
(3.27)
The singularity is at the points $`(M_{22}M_{33}\zeta _1^2,M_{11}M_{33},M_{11}M_{22},2B)=(0,0,0,0)`$. This corresponds to a $`\mathrm{C}\text{ }^{}`$ of singularities. Changing the variables, such that the singularity will pass through the origin, one has:
$`M_{22}`$ $`=`$ $`\zeta _1+x+y,`$ (3.28)
$`M_{33}`$ $`=`$ $`\zeta _1+xy.`$
Taking the lowest order in the polynomial (just for analyzing the points near the singularity), gives:
$$M_{11}(2x\zeta _1^2)=B^2.$$
(3.29)
That is a $`\text{/}\text{\_Z}_2`$ singularity. In order to obtain the field theory, let us give the following vev’s to the fields:
$$Z_1=0,Z_2=v\sigma _3,Z_3=v\sigma _3.$$
(3.30)
Plugging these values back into the F-flatness conditions, allows us to get the deformations:
$$\zeta _2=\zeta _3=0,\zeta _1=2v^2.$$
(3.31)
These vev’s break the group down to $`U(n/2)\times U(n/2)`$. The remaining spectrum is:
| | $`U(n/2)`$ | $`U(n/2)`$ |
| --- | --- | --- |
| $`\varphi _1`$ | adj | 1 |
| $`\varphi _2`$ | 1 | adj |
| $`A_{12}`$ | | $`\overline{\text{ }\text{ }\text{ }\text{ }\text{ }}`$ |
| $`B_{12}`$ | | $`\overline{\text{ }\text{ }\text{ }\text{ }\text{ }}`$ |
| $`A_{21}`$ | $`\overline{\text{ }\text{ }\text{ }\text{ }\text{ }}`$ | |
| $`B_{21}`$ | $`\overline{\text{ }\text{ }\text{ }\text{ }\text{ }}`$ | |
The $`A_{12}`$, $`B_{12}`$, $`A_{21}`$, $`B_{21}`$, $`\varphi _1`$ and $`\varphi _2`$ can be seen as the components of the $`Z_i`$ matrices:
$$Z_1=\left(\begin{array}{cc}B_{11}& B_{12}\\ B_{21}& B_{22}\end{array}\right),$$
(3.32)
$$Z_2=v\sigma _3+\left(\begin{array}{cc}(M_1+\varphi _1)/2& 2A_{12}\\ 2A_{21}& (M_2+\varphi _2)/2\end{array}\right),$$
(3.33)
$$Z_3=v\sigma _3+\left(\begin{array}{cc}(M_1\varphi _1)/2& 2A_{12}\\ 2A_{21}& (M_2\varphi _2)/2\end{array}\right).$$
(3.34)
The fields $`M_i`$ and $`B_{ii}`$ get masses. The relation between the non-diagonal components of the $`Z_2`$ and $`Z_3`$ fields comes from the degrees of freedom that are eaten by the Higgs mechanism. The superpotential for the massless fields is obtained by integrating out the massive states:
$$W=2Tr[(\varphi _1+\varphi _2)(A_{12}B_{21}+A_{21}B_{12})].$$
(3.35)
* Let us now take all the three $`\zeta _i`$ different from zero. For that, one can analyze how are the possible deformations of the $`\text{/}\text{\_Z}_2`$ case above (another possibility with the same result is taking the vev’s of the fields all proportional to $`\sigma _3`$). Let us take for simplicity $`U(1)\times U(1)`$ as the gauge group at the $`\text{/}\text{\_Z}_2`$ singularity. The deformations can be expressed in terms of the massless fields using the F-flatness conditions. There are four possible solutions:
$$\mathrm{\Delta }W=\pm 2v^2[4v+\frac{8}{v}A_{12}A_{21}+\frac{1}{4v}(\varphi _1^2+\varphi _2^2)],$$
(3.36)
$$\mathrm{\Delta }W=\pm 2v^2\frac{1}{4v}(\varphi _1^2\varphi _2^2).$$
(3.37)
The second one is the deformation considered by Klebanov and Witten that takes to the theory at the conifold. The superpotential gives masses to the two adjoints, leaving as massless spectrum:
| | $`U(n/2)`$ | $`U(n/2)`$ |
| --- | --- | --- |
| $`A_{12}`$ | | $`\overline{\text{ }\text{ }\text{ }\text{ }\text{ }}`$ |
| $`B_{12}`$ | | $`\overline{\text{ }\text{ }\text{ }\text{ }\text{ }}`$ |
| $`A_{21}`$ | $`\overline{\text{ }\text{ }\text{ }\text{ }\text{ }}`$ | |
| $`B_{21}`$ | $`\overline{\text{ }\text{ }\text{ }\text{ }\text{ }}`$ | |
## 4 Non-compact orientifold construction
Let us locate a set of $`D3`$-branes at an orbifold singularity of $`\mathrm{C}\text{ }^3/\mathrm{\Gamma }`$. The action of the orbifold group $`\mathrm{\Gamma }`$ permutes the branes in an appropriate way . We take an orientifold of the form $`\mathrm{\Omega }^{}=\mathrm{\Omega }(1)^{F_L}R_1R_2R_3`$ in order to preserve $`𝒩=1`$ supersymmetry in the effective theory in four dimensions (A models in ). The action of $`\mathrm{\Omega }`$ on the Chan-Paton matrices can be chosen to be either symmetric (A1 models) or antisymmetric (A2 models). In general these models will need the presence of $`D7`$-branes for consistency. This can be deduced from the tadpole cancellation conditions.
### 4.1 Closed string spectrum
As in the orbifold case, the closed string spectrum is continuous because the space $`\mathrm{C}\text{ }^3/\mathrm{\Gamma }`$ is non-compact. Again, the philosophy is to find the spectrum that would exist if the internal space were compactified. The $`\mathrm{\Omega }`$ parity is accompanied by a $`J`$ operation that relates states from one twisted sector with states from the inverse sector. In order to obtain the $`\mathrm{\Omega }J`$ invariant states one must combine the states from the $`g`$ and $`g^1`$ twisted sectors. Because of this, the counting of states is different compared to the orbifold case.
For sectors that are invariant under the $`J`$ permutation, i.e. the untwisted sector and the order-two twisted sectors, the spectrum in four dimensions is obtained by dimensionally reducing the fields of the type I string in ten dimensions. These fields are the metric $`g^{(10)}`$ and the dilaton $`\varphi ^{(10)}`$ in the NSNS sector and the 2-form $`C_2^{(10)}`$ in the RR sector. Their Lorentz indices have to be contracted with the differential forms of the untwisted and order-two twisted cohomology (see section 3.1). In $`𝒩=1`$ representations these sectors give: the gravity multiplet, a linear multiplet and $`(h^{1,1}+h^{2,1})_{\mathrm{untw}+\text{order-two}}`$ chiral multiplets.
One can split the remaining sectors into two types:
\- The sectors twisted by $`g`$ that give the same contribution to the cohomology as their inverse sectors, i.e. twisted by $`g^1`$ (these sectors always have fixed planes). Together they give $`h_g^{1,1}`$ linear multiplets, $`h_g^{2,1}`$ vectors and $`h_g^{1,1}+h_g^{2,1}`$ chiral multiplets .
\- The sectors twisted by $`g`$ that give a different contribution to the cohomology as their inverse sectors. Only one combination of the fields survives: $`h_g^{1,1}`$ linear multiplets and $`h_g^{2,1}`$ vectors. If the contribution from the $`g`$-twisted sector to $`h_g^{1,2}`$ and $`h_g^{2,1}`$ is not the same (that happens if the discrete torsion is different from $`\pm 1`$), one cannot match properly the states from one sector to the inverse sector. An orientifold model that contains such a sector is ill-defined. This consistency criterium agrees with the one we will find in the open string sector. There, it turns out that only real values of the discrete torsion are allowed.<sup>5</sup><sup>5</sup>5We thank Angel Uranga for pointing out to us that this can also be seen by noticing that $`\mathrm{\Omega }J`$ is not a symmetry of type IIB theory in the case of non-real discrete torsion.
In table 2, one can compare the $`\text{/}\text{\_Z}_2\times \text{/}\text{\_Z}_2`$ and $`\text{/}\text{\_Z}_3\times \text{/}\text{\_Z}_3`$ closed string spectrum. The first one admits the values $`ϵ=\pm 1`$ for the discrete torsion. Note that, although the cohomology is different in both cases, the closed string spectrum is the same for both values of the discrete torsion. The $`\text{/}\text{\_Z}_3\times \text{/}\text{\_Z}_3`$ model has $`H^2(\mathrm{\Gamma },U(1))=\text{/}\text{\_Z}_3`$, that is $`ϵ=e^{\pm 2\pi i/3}`$ or $`ϵ=1`$. Only the last value is allowed in the orientifold case.
Let us explain in detail why the closed string spectrum cannot be consistently defined if $`h_g^{2,1}h_g^{1,2}`$ for a sector twisted by $`g`$. The states corresponding to this sector can be obtained from the shift formalism.<sup>6</sup><sup>6</sup>6This will be explained in detail in . Take one sector $`g`$ with $`h_g^{1,2}0`$ and $`h_g^{2,1}=0`$. This gives $`h_g^{1,2}`$ helicity components $`+1`$ of some vectors and none with helicity $`1`$. In the inverse sector $`g^1`$ one gets $`h_{g^1}^{1,2}=0`$ and $`h_{g^1}^{2,1}=h_g^{1,2}`$. This leads to $`h_{g^1}^{2,1}`$ helicity components $`1`$ of some vectors. In type IIB theory this causes no problem, we get $`h_g^{1,2}`$ vectors. The problem arises if we want to impose the $`J`$ projection that relates the $`g`$-twisted sector with its inverse: One cannot match the states from these two sectors because they have opposite helicity.
This rules out several orientifold models: $`\text{/}\text{\_Z}_3\times \text{/}\text{\_Z}_3`$ with discrete torsion, $`\text{/}\text{\_Z}_4\times \text{/}\text{\_Z}_4`$ with discrete torsion $`\pm i`$, $`\text{/}\text{\_Z}_3\times \text{/}\text{\_Z}_6`$ with discrete torsion, etc. One can check that the orientifold projection is compatible with discrete torsion only for the values $`ϵ=+1`$ and $`ϵ=1`$. This can be seen from the formulae for $`h_g^{1,2}`$ and $`h_g^{2,1}`$ in the compact cases :
$`h_g^{1,2}`$ $`=`$ $`{\displaystyle \frac{1}{|\mathrm{\Gamma }|}}{\displaystyle \underset{h\mathrm{\Gamma }}{}}\beta _{g,h}\stackrel{~}{\chi }(g,h)e^{2\pi iv_h(g)},`$ (4.1)
$`h_g^{2,1}`$ $`=`$ $`{\displaystyle \frac{1}{|\mathrm{\Gamma }|}}{\displaystyle \underset{h\mathrm{\Gamma }}{}}\beta _{g,h}\stackrel{~}{\chi }(g,h)e^{2\pi iv_h(g)}.`$
Only if the discrete torsion is real, are $`h_g^{1,2}`$ and $`h_g^{2,1}`$ equal for any twist. We will see that this criterium in the closed string sector coincides with the one coming from the open string sector related to projective representations.
### 4.2 Tadpoles
As stated above, discrete torsion appears as a phase between different twisted sectors in closed string theories and as a phase between the matrices that represent the action of $`\mathrm{\Gamma }`$ on the Chan-Paton indices in the open string sector. In order to compute the tadpole contribution, one must sum over three different diagrams: the cylinder ($`𝒞`$), the Möbius strip ($``$) and the Klein bottle ($`𝒦`$). The first two diagrams contain a trace over the projective representation $`\gamma `$. Discrete torsion phases appear in the Möbius strip and in the Klein bottle. The three diagrams can be written as a sum over twisted sectors:
$`𝒞={\displaystyle \underset{g\mathrm{\Gamma }}{}}\stackrel{~}{𝒞}(g),={\displaystyle \underset{g\mathrm{\Gamma }}{}}\stackrel{~}{}(g),𝒦={\displaystyle \underset{g,h\mathrm{\Gamma }}{}}\stackrel{~}{ϵ}_h\beta _{h,g}\stackrel{~}{𝒦}_h(g),`$
$`\mathrm{with}\stackrel{~}{𝒞}(g)(Tr\gamma _g)^2\mathrm{and}\stackrel{~}{}(g)Tr(\gamma _{\mathrm{\Omega }g}^1\gamma _{\mathrm{\Omega }g}^{}).`$
The phases $`\beta _{h,g}`$ satisfy (2.2) and are related to the discrete torsion. The additional factor $`\stackrel{~}{ϵ}_h=\pm 1`$ appears because of the possibility of having vector structure or not in each of the two factors of $`\text{/}\text{\_Z}_N\times \text{/}\text{\_Z}_M`$. As $`J`$ relates the orbifold elements $`h`$ and $`h^1`$, the only non-vanishing contributions to the Klein bottle are (see appendix of ) $`\stackrel{~}{𝒦}_e(g)`$ and $`\stackrel{~}{𝒦}_{h^{(2)}}(g)`$, where $`e`$ is the neutral element of $`\mathrm{\Gamma }`$ and $`h^{(2)}`$ is of order two, i.e. $`(h^{(2)})^2=e`$. The first of these two contributions is not modified by discrete torsion, because (2.2) implies $`\beta _{e,g}=1`$. The second contribution depends on discrete torsion, as explained in appendix B.
In general, models of this type require the presence of $`D7`$-branes. The tadpole calculation is detailed in appendix B. To give the result of this calculation, we first introduce some notation. It is convenient to denote the group elements of $`\mathrm{\Gamma }=\text{/}\text{\_Z}_N\times \text{/}\text{\_Z}_M`$ by $`\overline{k}=(n,m)`$ as in eq. (3.1). One has to distinguish even twists $`\overline{k}`$, i.e. there is a $`\overline{k}^{}\mathrm{\Gamma }`$, such that $`\overline{k}=2\overline{k}^{}`$, from odd twists that cannot be written in this form. If $`N`$ is odd and $`M`$ is even, there exists one order-two element $`\overline{k}_1`$. If $`N`$ and $`M`$ are both even there are three order-two elements $`\overline{k}_1`$, $`\overline{k}_2`$, $`\overline{k}_3`$, where the index denotes the complex plane that is fixed by this element (i.e. $`\overline{k}_1=(0,M/2)`$, etc). As will be explained below, the generalization of the notion of vector structure to $`\text{/}\text{\_Z}_N\times \text{/}\text{\_Z}_M`$ orientifolds leads to four different boundary conditions on the $`\gamma `$ matrices: $`\gamma _{g_1}^N=\pm 1\mathrm{I}`$, $`\gamma _{g_2}^M=\pm 1\mathrm{I}`$. We define the numbers $`\mu _i=\pm 1`$ (the index $`i`$ refers to the complex plane that is fixed by the corresponding group element) by
$$\gamma _{g_1}^N=\mu _3\mathrm{\hspace{0.17em}1}\mathrm{I},\gamma _{g_2}^M=\mu _1\mathrm{\hspace{0.17em}1}\mathrm{I},(\gamma _{g_1}\gamma _{g_2})^{\mathrm{lcm}(N,M)}=\mu _21\mathrm{I}.$$
(4.3)
Of course, only two of the $`\mu _i`$ are independent. One can show that
$$\mu _2=\{\begin{array}{cc}ϵ^{MN/4}\mu _3\mu _1\hfill & \mathrm{if}\frac{M}{\mathrm{gcd}(N,M)}\mathrm{and}\frac{N}{\mathrm{gcd}(N,M)}\mathrm{odd}\hfill \\ \mu _1\hfill & \mathrm{if}\frac{M}{\mathrm{gcd}(N,M)}\mathrm{even}\hfill \\ \mu _3\hfill & \mathrm{if}\frac{N}{\mathrm{gcd}(N,M)}\mathrm{even}\hfill \end{array},$$
(4.4)
where $`ϵ`$ is the discrete torsion parameter of (2.3) and we used that $`ϵ`$ only takes the values $`\pm 1`$. Due to the fact that the $`\gamma `$-matrices form a projective representation, there may appear a phase $`\delta _{\overline{k}}`$ in the Möbius strip:
$$Tr(\gamma _{\mathrm{\Omega }\overline{k}}^1\gamma _{\mathrm{\Omega }\overline{k}}^{})=\delta _{\overline{k}}Tr(\gamma _{\overline{k}}^2).$$
(4.5)
Finally, we define
$$s_i=\mathrm{sin}(\pi \overline{k}\overline{v}_i),\stackrel{~}{s}_i=\mathrm{sin}(2\pi \overline{k}\overline{v}_i),c_i=\mathrm{cos}(\pi \overline{k}\overline{v}_i).$$
(4.6)
In this notation the tadpole conditions for a general $`\text{/}\text{\_Z}_N\times \text{/}\text{\_Z}_M`$ orientifold read:
* If $`N`$ and $`M`$ are both odd:
$$8\stackrel{~}{s}_1\stackrel{~}{s}_2\stackrel{~}{s}_3Tr\gamma _{2\overline{k},3}+\underset{i=1}{\overset{3}{}}2\stackrel{~}{s}_iTr\gamma _{2\overline{k},7_i}=\delta _{\overline{k}}32s_1s_2s_3.$$
(4.7)
* If $`N`$ is odd and $`M`$ is even:
+ for odd $`\overline{k}`$:
$$8s_1s_2s_3Tr\gamma _{\overline{k},3}+\underset{i=1}{\overset{3}{}}2s_iTr\gamma _{\overline{k},7_i}=0,$$
(4.8)
+ for even twists $`2\overline{k}`$:
$$8\stackrel{~}{s}_1\stackrel{~}{s}_2\stackrel{~}{s}_3Tr\gamma _{2\overline{k},3}+\underset{i=1}{\overset{3}{}}2\stackrel{~}{s}_iTr\gamma _{2\overline{k},7_i}=\delta _{\overline{k}}32(s_1s_2s_3\stackrel{~}{ϵ}_1s_1c_2c_3),$$
(4.9)
where $`\stackrel{~}{ϵ}_1=\stackrel{~}{ϵ}_{\overline{k}_1}`$.
As explained in appendix B, in order to factorize the amplitudes in the appropriate way, one must impose some restrictions on the ‘vector structures’. In the present case, one needs<sup>7</sup><sup>7</sup>7Strictly speaking, these relations only follow from tadpole factorization if $`M>2`$. However, these conditions are still valid in the case $`M=2`$. This can be seen by similar arguments as the ones used in . The same applies to the relations (c)) below.:
$$\mu _{1}^{}{}_{}{}^{3}=\mu _{1}^{}{}_{}{}^{7_1}=\mu _{1}^{}{}_{}{}^{7_2}=\mu _{1}^{}{}_{}{}^{7_3}=\stackrel{~}{ϵ}_1.$$
(4.10)
* If $`N`$ and $`M`$ are both even:
Only in this case discrete torsion $`ϵ=1`$ is possible. The tadpole conditions are:
+ for odd $`\overline{k}`$:
$$8s_1s_2s_3Tr\gamma _{\overline{k},3}+\underset{i=1}{\overset{3}{}}2s_iTr\gamma _{\overline{k},7_i}=0,$$
(4.11)
+ for even twists $`2\overline{k}`$:
$$8\stackrel{~}{s}_1\stackrel{~}{s}_2\stackrel{~}{s}_3Tr\gamma _{2\overline{k},3}+\underset{i=1}{\overset{3}{}}2\stackrel{~}{s}_iTr\gamma _{2\overline{k},7_i}=ϵ^{k_1k_2}\delta _{\overline{k}}32(s_1s_2s_3\underset{ijk}{}\stackrel{~}{ϵ}_i\beta _is_ic_jc_k),$$
(4.12)
where $`\stackrel{~}{ϵ}_i=\stackrel{~}{ϵ}_{\overline{k}_i}`$, $`\beta _i=\beta _{\overline{k}_i,\overline{k}}`$.
As in the previous case, there are some conditions on the ‘vector structures’:
$`\mu _{1}^{}{}_{}{}^{3}=\mu _{1}^{}{}_{}{}^{7_1}=\mu _{1}^{}{}_{}{}^{7_2}=\mu _{1}^{}{}_{}{}^{7_3}=\stackrel{~}{ϵ}_1,`$
$`\mu _{3}^{}{}_{}{}^{3}=\mu _{3}^{}{}_{}{}^{7_1}=\mu _{3}^{}{}_{}{}^{7_2}=\mu _{3}^{}{}_{}{}^{7_3}=\stackrel{~}{ϵ}_3.`$ (4.13)
In addition the $`\stackrel{~}{ϵ}_i`$ are related by
$$\stackrel{~}{ϵ}_1\stackrel{~}{ϵ}_2\stackrel{~}{ϵ}_3=ϵ^{MN/4}.$$
(4.14)
Let us discuss some examples:
* $`\text{/}\text{\_Z}_2\times \text{/}\text{\_Z}_2`$: All tadpole conditions vanish. As in the orbifold case, the contribution from the untwisted sector vanishes because the space $`\mathrm{C}\text{ }^3/\mathrm{\Gamma }`$ transverse to the $`D3`$-branes is non-compact and the tadpoles depend on the inverse volume. The twisted sectors $`(1,0)`$, $`(0,1)`$ and $`(1,1)`$ each leave one complex plane fixed (the third, the first and the second respectively). The sum over the windings along the fixed direction leads to a dependence on the inverse volume of the fixed plane. In particular, no $`D7`$-branes are needed for this model to be consistent.
* $`\text{/}\text{\_Z}_3\times \text{/}\text{\_Z}_3`$: There is one twisted sector that does not fix any plane but has only a fixed point at the origin, namely the sector twisted by $`g=(1,2)`$ or by its inverse $`g^1=(2,1)`$. Tadpole cancellation implies :
$$Tr\gamma _{(1,2),3}+\frac{1}{3}(Tr\gamma _{(1,2),7_1}+Tr\gamma _{(1,2),7_2}Tr\gamma _{(1,2),7_3})=4\delta _{(2,1)}.$$
(4.15)
But this condition cannot be satisfied if there is discrete torsion. The only allowed values of the discrete torsion parameter different from one are $`e^{\pm 2\pi i/3}`$. For these values there is a unique projective representation (see appendix A) with the following character:
$$Tr\gamma _g=0ge.$$
(4.16)
($`Tr\gamma _e`$ is related to the total number of branes.) Thus, we see that the above condition can never be satisfied for this projective representation. This inconsistency agrees with the one we found in the closed string sector of the same model.
A similar analysis can be done for other orientifolds, with the result that no solution can be found for certain values of the discrete torsion. In $`\text{/}\text{\_Z}_N\times \text{/}\text{\_Z}_N`$ models , with discrete torsion $`ϵ=e^{2\pi im/N}`$ and $`m`$ and $`N`$ coprime, there is only one irreducible representation<sup>8</sup><sup>8</sup>8In general there are $`N^2/s^2`$ different irreducible representations. If, however, $`\mathrm{gcd}(N,m)=1`$, then the minimal positive integer $`s`$, such that $`ϵ^s=1`$, is $`s=N`$. of the discrete group with this torsion as a factor system. This means that the caracters of this representation vanish for all group elements except for the unit element, i.e. (4.16) is valid for all of these models. Following the same argument as in the $`\text{/}\text{\_Z}_3\times \text{/}\text{\_Z}_3`$ case, one finds that all $`\text{/}\text{\_Z}_N\times \text{/}\text{\_Z}_N`$, $`N>2`$, orientifolds with minimal discrete torsion (i.e. $`s=N`$) are inconsistent. Again, this result agrees with the restrictions found in the closed string sector.
* $`\text{/}\text{\_Z}_2\times \text{/}\text{\_Z}_4`$: The tadpole conditions are:
$$\begin{array}{ccc}\overline{k}\hfill & \overline{k}\overline{v}\hfill & \text{tadpole condition}\hfill \\ & & \\ (0,1)\hfill & \frac{1}{4}(0,1,1)\hfill & Tr\gamma _{(0,1),7_2}Tr\gamma _{(0,1),7_3}=0\hfill \\ (0,3)\hfill & \frac{3}{4}(0,1,1)\hfill & Tr\gamma _{(0,3),7_2}Tr\gamma _{(0,3),7_3}=0\hfill \\ (1,0)\hfill & \frac{1}{2}(1,1,0)\hfill & Tr\gamma _{(1,0),7_1}Tr\gamma _{(1,0),7_2}=0\hfill \\ (1,1)\hfill & \frac{1}{4}(2,1,1)\hfill & 4Tr\gamma _{(1,1),3}+2Tr\gamma _{(1,1),7_1}\sqrt{2}Tr\gamma _{(1,1),7_2}\sqrt{2}Tr\gamma _{(1,1),7_3}=0\hfill \\ (1,2)\hfill & \frac{1}{2}(1,0,1)\hfill & Tr\gamma _{(1,2),7_1}Tr\gamma _{(1,2),7_3}=0\hfill \\ (1,3)\hfill & \frac{1}{4}(2,1,3)\hfill & 4Tr\gamma _{(1,3),3}+2Tr\gamma _{(1,3),7_1}+\sqrt{2}Tr\gamma _{(1,3),7_2}\sqrt{2}Tr\gamma _{(1,3),7_3}=0\hfill \\ (0,2)\hfill & \frac{1}{2}(0,1,1)\hfill & 2Tr\gamma _{(0,2),7_2}2Tr\gamma _{(0,2),7_3}=16\delta _{(0,1)}(\stackrel{~}{ϵ}_2\beta _2\stackrel{~}{ϵ}_3\beta _3)\hfill \end{array}$$
There are two cases:
(i) If $`\stackrel{~}{ϵ}_2\beta _2=\stackrel{~}{ϵ}_3\beta _3`$, then the Klein bottle contribution in the $`(0,2)`$ sector vanishes. This means that no $`D7`$-branes are needed to cancel the tadpoles. This agrees with anomaly cancellation in the 33 sector. Using $`\beta _2=\beta _{(1,2),(0,1)}=ϵ`$, $`\beta _3=\beta _{(1,0),(0,1)}=ϵ`$ and eqs. (c)), (4.14), we find that $`\stackrel{~}{ϵ}_2\beta _2=\stackrel{~}{ϵ}_3\beta _3`$ implies $`\mu _{1}^{}{}_{}{}^{3}\mu _{3}^{}{}_{}{}^{3}=\mu _{3}^{}{}_{}{}^{3}`$. Therefore this case is characterized by $`\mu _{1}^{}{}_{}{}^{3}=+1`$.
(ii) If $`\stackrel{~}{ϵ}_2\beta _2=\stackrel{~}{ϵ}_3\beta _3`$, then $`D7`$-branes are present. In the case with discrete torsion $`ϵ=1`$, a minimal choice to satisfy the tadpole conditions consists in a set of $`D7_2`$-branes. Then all the traces vanish, except for the $`(0,2)`$ sector, as can be seen from (A.19). With the minimal choice, one has
$$Tr\gamma _{(0,2),7_2}=16\delta _{(0,1)}\stackrel{~}{ϵ}_3\beta _3=16\delta _{(0,1)}\mu _{3}^{}{}_{}{}^{3},$$
(4.17)
where we used (c)) and $`\beta _3=\beta _{(1,0),(0,1)}=1`$. The relation $`\stackrel{~}{ϵ}_2\beta _2=\stackrel{~}{ϵ}_3\beta _3`$ implies $`\mu _{1}^{}{}_{}{}^{3}\mu _{3}^{}{}_{}{}^{3}=\mu _{3}^{}{}_{}{}^{3}`$ and thus $`\mu _{1}^{}{}_{}{}^{3}=1`$. From (c)) we then find that this case is characterized by $`\mu _{1}^{}{}_{}{}^{3}=\mu _{1}^{}{}_{}{}^{7_2}=1`$.
### 4.3 open string spectrum
An element $`g`$ of the orbifold group $`\mathrm{\Gamma }`$ and the world-sheet parity $`\mathrm{\Omega }`$ act on the Chan-Paton matrices $`\lambda `$ as
$$g:\lambda \gamma _g\lambda \gamma _g^1,\mathrm{\Omega }:\lambda \gamma _\mathrm{\Omega }\lambda ^{}\gamma _\mathrm{\Omega }^1.$$
(4.18)
The open string spectrum, i.e. gauge fields $`\lambda ^{(0)}`$ and matter fields $`\lambda ^{(i)}`$, is obtained by taking the solutions to the projections:
$$\gamma _g^1\lambda \gamma _g=r(g)\lambda ,\gamma _\mathrm{\Omega }^1\lambda ^{}\gamma _\mathrm{\Omega }=r(\mathrm{\Omega })\lambda ,$$
(4.19)
where $`r(g)`$ (resp. $`r(\mathrm{\Omega })`$) is the matrix that represents the action of $`g`$ (resp. $`\mathrm{\Omega }`$) on $`\lambda =(\lambda ^{(0)},\lambda ^{(1)},\lambda ^{(2)},\lambda ^{(3)})`$.
The relation $`\mathrm{\Omega }^2=e`$ gives a restriction on the matrix $`\gamma _\mathrm{\Omega }`$:
$$\gamma _\mathrm{\Omega }(\gamma _\mathrm{\Omega }^1)^{}=c\mathrm{\hspace{0.17em}1}\mathrm{I},$$
(4.20)
where the constant $`c`$ can only take the values $`\pm 1`$, i.e. $`\gamma _\mathrm{\Omega }`$ is either symmetric or antisymmetric. A further condition follows from $`\mathrm{\Omega }g\mathrm{\Omega }=g`$, which, using (4.18), translates to:
$$\gamma _\mathrm{\Omega }(\gamma _g^1)^{}(\gamma _\mathrm{\Omega }^1)^{}=\delta _g\gamma _g,$$
(4.21)
where the phase $`\delta _g`$ appears because the representation is only projective. More precisely, in the notation of appendix A, $`\delta _g=c\beta _{\mathrm{\Omega },g}`$. This can be seen as follows. Because of the special action of $`\mathrm{\Omega }`$ on $`\lambda `$, eq. (4.18), the matrices representing $`\mathrm{\Omega }g`$ and $`g\mathrm{\Omega }`$ are given by
$$\gamma _{\mathrm{\Omega }g}=\alpha _{\mathrm{\Omega },g}^1\gamma _\mathrm{\Omega }(\gamma _g^1)^{},\gamma _{g\mathrm{\Omega }}=\alpha _{g,\mathrm{\Omega }}^1\gamma _g\gamma _\mathrm{\Omega }.$$
(4.22)
From $`\mathrm{\Omega }g=g\mathrm{\Omega }`$ and (4.20), we find $`\delta _g=c\alpha _{\mathrm{\Omega },g}\alpha _{g,\mathrm{\Omega }}^1`$. It is easy to see that the phase $`\delta _g`$ coincides with the one defined in (4.5).
As we are interested in unitary projective representations, (4.21) can be transformed into
$$\gamma _\mathrm{\Omega }\gamma _g^{}\gamma _\mathrm{\Omega }^1=c\delta _g\gamma _g,$$
(4.23)
where we used $`\gamma _\mathrm{\Omega }^{}=c\gamma _\mathrm{\Omega }`$. The world-sheet parity $`\mathrm{\Omega }`$ relates one representation with its complex conjugate. This means that the matrices $`\gamma _g`$ of the orientifold form a (pseudo-)real projective representation of the orbifold group. As the complex projective representations are classified by $`H^2(\mathrm{\Gamma },\mathrm{C}\text{ }^{})=H^2(\mathrm{\Gamma },U(1))`$, the real projective representations are classified by $`H^2(\mathrm{\Gamma },\mathrm{IR}^{})=H^2(\mathrm{\Gamma },\text{/}\text{\_Z}_2)`$. As a consequence, the discrete torsion paramter $`ϵ`$ can only take the real values $`\pm 1`$.
Let us make one further remark concerning the phases $`\beta _{\mathrm{\Omega },g}`$. They do not satisfy (A.5) but rather $`\beta _{\mathrm{\Omega },gh}=(\alpha _{g,h})^2\beta _{\mathrm{\Omega },g}\beta _{\mathrm{\Omega },h}`$. Moreover, one can always find an equivalent set of $`\gamma `$-matrices with $`\beta _{\mathrm{\Omega },g}=1g\mathrm{\Gamma }`$ by defining $`\widehat{\gamma }_g=\sqrt{\beta _{\mathrm{\Omega },g}}\gamma _g`$. If we choose the factor system $`\alpha _{g,h}`$ of the orbifold group $`\mathrm{\Gamma }`$ as in (A.6), then $`(\alpha _{g,h})^2=1g,h`$ and therefore the $`\widehat{\gamma }_g`$ have the same factor system as the $`\gamma _g`$: $`\widehat{\alpha }_{g,h}=\sqrt{\beta _{\mathrm{\Omega },g}\beta _{\mathrm{\Omega },h}\beta _{\mathrm{\Omega },gh}^1}\alpha _{g,h}=\alpha _{g,h}`$. Thus, to fix the factor system of the orientifold completely, one has to add one relation to (A.6):
$`\alpha _{g_1^a,g_2^b}=\alpha _{g_1,g_1^a}=\alpha _{g_2,g_2^b}`$ $`=`$ $`1,a=1,\mathrm{},N,b=1,\mathrm{},M,`$
$`\delta _g`$ $`=`$ $`cg\mathrm{\Gamma }.`$ (4.24)
In the following, we will mostly leave $`\alpha _{g,h}`$, $`\delta _g`$ arbitrary. But to perform calculations, we choose the factor system (4.3).
The reality condition (4.23) has an interesting consequence concerning the classification of $`\text{/}\text{\_Z}_N`$ orientifolds. Note, that for a complex projective representation of $`\text{/}\text{\_Z}_N`$ the relation
$$\gamma _g^N=\stackrel{~}{c}\mathrm{\hspace{0.17em}1}\mathrm{I},\stackrel{~}{c}\mathrm{C}\text{ }^{},$$
(4.25)
can always be brought to the form
$$\widehat{\gamma }_g^N=1\mathrm{I},$$
(4.26)
defining $`\widehat{\gamma }_g=\stackrel{~}{c}^{1/N}\gamma _g`$. For real projective representations, i.e. $`\stackrel{~}{c}\mathrm{IR}^{}`$, this is only possible if $`N`$ is odd or if $`\stackrel{~}{c}>0`$. If $`N`$ is even, two inequivalent projective representations arise:
$$\gamma _g^N=\pm 1\mathrm{I}.$$
(4.27)
These are the cases with and without vector structure .
If $`\mathrm{\Gamma }=\text{/}\text{\_Z}_N\times \text{/}\text{\_Z}_M`$ (with $`N,M`$ both even), then four cases have to be distinguished: $`\gamma _{g_1}^N=\pm 1\mathrm{I}`$, $`\gamma _{g_2}^M=\pm 1\mathrm{I}`$. Only if $`\gamma _{g_1}^N=\gamma _{g_2}^M=1\mathrm{I}`$ and $`\gamma _{g_1}\gamma _{g_2}=\gamma _{g_2}\gamma _{g_1}`$, does the gauge bundle of the orientifold have vector structure in the sense of . Thus, orientifolds with discrete torsion can never have vector structure. However, the different boundary conditions for the gauge bundle of orientifolds with discrete torsion — let us denote them by $`(++)`$, $`(+)`$, $`(+)`$ and $`()`$ — lead, in general, to non-equivalent models. In particular, the gauge group and the spectrum are not the same in all four cases.
As explained in appendix A, there are two types of real projective representations: irreducible real projective representations and combinations of pairs of conjugate irreducible complex representations. Note, that if $`\stackrel{~}{\gamma }`$ is a complex projective representation of the discrete group $`\mathrm{\Gamma }`$, so is $`\stackrel{~}{\gamma }^{}=(\stackrel{~}{\gamma }^{})^1`$. The orientifold projection tells us that one has to take these two representations together, whenever $`\stackrel{~}{\gamma }\stackrel{~}{\gamma }^{}`$. In terms of matrices this means:
$$\gamma _g=\left(\begin{array}{cc}\stackrel{~}{\gamma }_g& 0\\ 0& c\delta _g^1(\stackrel{~}{\gamma }_g^{})^1\end{array}\right)g\mathrm{\Gamma },$$
(4.28)
where $`c=\pm 1`$ and $`\delta _g`$ is a phase. It can be verified that $`\gamma _g`$ satisfies (4.21) if $`\gamma _\mathrm{\Omega }`$ is of the form
$$\gamma (\mathrm{\Omega })=\left(\begin{array}{cc}0& 1\\ c& 0\end{array}\right)1\mathrm{I}_n,$$
(4.29)
where $`n`$ is the dimension of $`\stackrel{~}{\gamma }`$. The fact that $`\gamma _\mathrm{\Omega }`$ can always be chosen to be of this form is a consequence of (4.20).
Now, let us see the restrictions that arise from the multiplication of two group elements $`g,h\mathrm{\Gamma }`$:
$$\gamma _g\gamma _h=\left(\begin{array}{cc}\stackrel{~}{\gamma }_g\stackrel{~}{\gamma }_h& 0\\ 0& \delta _g^1\delta _h^1(\stackrel{~}{\gamma }_g^{})^1(\stackrel{~}{\gamma }_h^{})^1\end{array}\right)$$
(4.30)
From the relation $`\gamma _g\gamma _h=\alpha _{g,h}\gamma _{gh}`$, eq. (3.6), one finds
$$\gamma _g\gamma _h=\alpha _{g,h}\left(\begin{array}{cc}\stackrel{~}{\gamma }_{gh}& 0\\ 0& \delta _g^1\delta _h^1\alpha _{g,h}^2(\stackrel{~}{\gamma }_{gh}^{})^1\end{array}\right)=\alpha _{g,h}\gamma _{gh},$$
(4.31)
which gives the relation between the discrete torsion $`\alpha _{g,h}`$ and the orientifold phase $`\delta _g`$:
$$\delta _g\delta _h\alpha _{g,h}^2=c\delta _{gh}.$$
(4.32)
This condition comes from the fact that two projective representation can be related if they belong to the same equivalence class of factor systems (see appendix A or ). The projective representation $`\stackrel{~}{\gamma }_g`$ has the factor system $`\alpha _{g,h}`$ and $`(\stackrel{~}{\gamma }_g^{})^1`$ the inverse one, $`\alpha _{g,h}^1`$. They are in the same equivalence class if there exist numbers $`\rho _g`$, such that $`(\rho _g\rho _h/\rho _{gh})\alpha _{g,h}=\alpha _{g,h}^1`$, which is just the above condition.
Interchanging $`g`$ and $`h`$ in (4.32), one finds $`(\delta _h\delta _g/c\delta _{hg})\alpha _{h,g}=\alpha _{h,g}^1`$. The quotient of these two relations gives the following consistency condition: $`(\alpha _{h,g}\alpha _{g,h}^1)^2=1`$. This tells us that the only values of discrete torsion compatible with the orientifold projection are
$$\beta _{g,h}=\alpha _{g,h}\alpha _{h,g}^1=\pm 1.$$
(4.33)
This result coincides with the consistency condition that we found in the closed string sector of orientifolds with discrete torsion.
According to the results of appendix A, the matrix $`\gamma _g`$ (whith $`g=g_1^ag_2^b`$) of a general real projective representation with discrete torsion $`ϵ=1`$ and with $`\eta _{1/2}=0,1`$, such that $`(\gamma _{g_1})^N=(1)^{\eta _1}1\mathrm{I}`$, $`(\gamma _{g_2})^M=(1)^{\eta _2}1\mathrm{I}`$, is of the form
$$\underset{k,l}{}\left((\omega _{2N}^{2k+\eta _1}\gamma _{g_1})^a(\omega _{2M}^{2l+\eta _2}\gamma _{g_2})^b(\omega _{2N}^{2N2k\eta _1}\gamma _{g_1})^a(\omega _{2M}^{2M2l\eta _2}\gamma _{g_2})^b\right)1\mathrm{I}_{n_{kl}},$$
(4.34)
where
$$\gamma _{g_1}=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right),\gamma _{g_2}=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right),$$
(4.35)
as given in (A.11), and $`k=0,\mathrm{},N/21`$, $`l=0,\mathrm{},M/21`$. Here, we chose the factor system (A.6), which implies $`(\gamma _{g_1})^a(\gamma _{g_2})^b=\gamma _{g_1^ag_2^b}`$.
Let us now analyze some examples. Start with the $`\text{/}\text{\_Z}_2\times \text{/}\text{\_Z}_2`$ case with $`(++)`$ boundary condition, i.e. $`\gamma _{g_1}^2=\gamma _{g_2}^2=1\mathrm{I}`$. Discrete torsion is possible for this model, because $`H^2(\text{/}\text{\_Z}_2\times \text{/}\text{\_Z}_2,\text{/}\text{\_Z}_2)=\text{/}\text{\_Z}_2`$. The case without discrete torsion has four irreducible representations, all of them are one dimensional and real. The case with discrete torsion $`1`$ has a unique irreducible representation, that can be taken to be real. A general representation, with $`(++)`$ boundary condition, is of the form
$`\gamma _e`$ $`=`$ $`1\mathrm{I}_21\mathrm{I}_n,`$ (4.36)
$`\gamma _{g_1}`$ $`=`$ $`\sigma _31\mathrm{I}_n,`$
$`\gamma _{g_2}`$ $`=`$ $`\sigma _11\mathrm{I}_n,`$
$`\gamma _{g_1g_2}`$ $`=`$ $`i\sigma _21\mathrm{I}_n,`$
where $`n`$ is an arbitrary parameter. If the matrix $`\gamma _\mathrm{\Omega }`$ is symmetric, then it can be taken to be of the form
$$\gamma _\mathrm{\Omega }=1\mathrm{I}_21\mathrm{I}_n.$$
(4.37)
If it is antisymmetric, we restrict ourselves to the case of even $`n`$ and take
$$\gamma _\mathrm{\Omega }=1\mathrm{I}_2\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)1\mathrm{I}_{n/2}.$$
(4.38)
Using (4.22), with $`\alpha _{\mathrm{\Omega },g}=\alpha _{g,\mathrm{\Omega }}=1`$, it is easy to check that all consistency conditions from the multiplication of group elements (like $`g_1^2=g_2^2=e`$, $`\mathrm{\Omega }g=g\mathrm{\Omega }`$, etc) are satisfied for this choice of matrices. As we mentioned above, this choice is unique up to equivalence.
The spectrum can be determined by taking the matrices $`\lambda ^{(0)},\lambda ^{(i)}`$ of the orbifold case (3.10),
$$\lambda ^{(0)}=1\mathrm{I}_2X,\lambda ^{(i)}=\sigma _iZ_i,$$
(4.39)
and restricting them, such that also the second condition in (4.19) is satisfied. If $`\gamma _\mathrm{\Omega }`$ is symmetric, then $`X`$, $`Z_1`$, $`Z_3`$ are antisymmetric and $`Z_2`$ is symmetric. This corresponds to gauge group $`SO(n)`$ with two adjoint fields, a traceless symmetric tensor and a singlet. If $`\gamma _\mathrm{\Omega }`$ is antisymmetric, we find<sup>9</sup><sup>9</sup>9In our notation, $`USp(2k)`$ is the unitary symplectic group of rank $`k`$. $`USp(n)`$ with two adjoint fields and an antisymmetric tensor. The superpotential is the one of the orbifold with the above representations:
$$W=Tr(Z_1Z_2Z_3+Z_2Z_1Z_3)$$
(4.40)
The solution to the case with $`()`$ boundary condition is essentially the $`\text{/}\text{\_Z}_2\times \text{/}\text{\_Z}_2`$ model of without $`D5`$-branes. In our conventions the $`\gamma `$-matrices are given by
$`\gamma _e`$ $`=`$ $`1\mathrm{I}_21\mathrm{I}_{2n},`$ (4.41)
$`\gamma _{g_1}`$ $`=`$ $`(i\sigma _1(i\sigma _1))1\mathrm{I}_n,`$
$`\gamma _{g_2}`$ $`=`$ $`(i\sigma _3(i\sigma _3))1\mathrm{I}_n,`$
$`\gamma _{g_1g_2}`$ $`=`$ $`i\sigma _21\mathrm{I}_{2n},`$
where we exchanged $`\gamma _{g_1}`$ and $`\gamma _{g_2}`$ with respect to (4.34) for later convenience. If $`\gamma _\mathrm{\Omega }`$ is symmetric, the spectrum consists of $`USp(2n)`$ gauge fields and three matter fields in the antisymmetric tensor representation. This is exactly the spectrum of the 99 sector of the model discussed in , if $`n=8`$. If $`\gamma _\mathrm{\Omega }`$ is antisymmetric, we find $`SO(2n)`$ with three symmetric tensors.
A similar analysis can be done for the boundary conditions $`(+)`$ and $`(+)`$. One finds in both cases the same spectrum as in the $`(++)`$ case.
To determine the spectrum of the $`\text{/}\text{\_Z}_2\times \text{/}\text{\_Z}_4`$ orientifold, it is more convenient to use the shift formalism, as explained in appendix C. We take $`\gamma _\mathrm{\Omega }`$ symmetric and boundary conditions $`(+)`$. Let $`g_1`$ resp. $`g_2`$ be the generators of $`\text{/}\text{\_Z}_2`$ resp. $`\text{/}\text{\_Z}_4`$. We choose a basis where $`\gamma _{g_2}`$ is diagonal. According to appendix C, $`g_2`$ is represented by the shift
$$V_{g_2}=\frac{1}{4}(0^{n_0},2^{n_0},1^{n_1},3^{n_1})$$
and $`g_1`$ is represented by the permutation, acting on the roots $`\rho `$ of the $`SO(2(n_0+n_1))`$ lattice,
$$\mathrm{\Pi }_{g_1}:(\stackrel{~}{\rho }_1^{(n_0)},\stackrel{~}{\rho }_2^{(n_0)},\stackrel{~}{\rho }_1^{(n_1)},\stackrel{~}{\rho }_2^{(n_1)})(\stackrel{~}{\rho }_2^{(n_0)},\stackrel{~}{\rho }_1^{(n_0)},\stackrel{~}{\rho }_2^{(n_1)},\stackrel{~}{\rho }_1^{(n_1)}),$$
where $`\stackrel{~}{\rho }_i^{(n_0)}=(\rho _{(i1)n_0+1},\mathrm{},\rho _{in_0})`$ and $`\stackrel{~}{\rho }_i^{(n_1)}=(\rho _{2n_0+(i1)n_1+1},\mathrm{},\rho _{2n_0+in_1})`$. As explained in appendix C no shift associated to $`g_1`$ is needed, whereas in the $`(++)`$ case we would need an additional shift $`V_{g_1}=\frac{1}{4}(1^{\stackrel{~}{n}_0})`$. (Note that $`g_1`$ and $`g_2`$ are exchanged with respect to the notation in appendix C). The gauge group is determined by finding all linear combinations of roots $`\rho `$ that satisfy $`\rho V_{g_2}=0`$ and $`\mathrm{\Pi }_{g_1}(\rho )=\rho `$. They are of the form
$`(\underset{¯}{+,+,0^{n_02}},0^{n_0},0^{2n_1})+(0^{n_0},\underset{¯}{+,+,0^{n_02}},0^{2n_1}),`$
$`(\underset{¯}{,,0^{n_02}},0^{n_0},0^{2n_1})+(0^{n_0},\underset{¯}{,,0^{n_02}},0^{2n_1}),`$
$`(\underset{¯}{+,,0^{n_02}},0^{n_0},0^{2n_1})+(0^{n_0},\underset{¯}{+,,0^{n_02}},0^{2n_1}),`$
$`(0^{2n_0},\underset{¯}{+,0^{n_11}},\underset{¯}{0^{n_11},+})+(0^{2n_0},\underset{¯}{0^{n_11},+},\underset{¯}{+,0^{n_11}}),`$
$`(0^{2n_0},\underset{¯}{,0^{n_11}},\underset{¯}{0^{n_11},})+(0^{2n_0},\underset{¯}{0^{n_11},},\underset{¯}{,0^{n_11}}),`$
$`(0^{2n_0},\underset{¯}{+,,0^{n_12}},0^{n_1})+(0^{2n_0},0^{n_1},\underset{¯}{+,,0^{n_12}}).`$
Underlining means that one has to take all correlated permutations of the two terms in each line, such that they are invariant under $`\mathrm{\Pi }_{g_1}`$. The first three lines give $`n_0(2n_01)n_0`$ roots, which, together with $`n_0`$ Cartan generators, form the gauge group $`SO(2n_0)`$. The last three lines give $`n_1(2n_1+1)n_1`$ roots, which, together with $`n_1`$ Cartan generators, form the gauge group $`USp(2n_1)`$. The matter fields from the first complex plane are obtained from the roots that satisfy $`\rho V_{g_2}=0`$ and $`\mathrm{\Pi }_{g_1}(\rho )=\rho `$. These are just the antisymmetric combinations of the roots above that formed the gauge group. Thus, we find an adjoint field of $`SO(2n_0)`$ and an antisymmetric tensor of $`USp(2n_1)`$. For the second complex plane, one has the condition $`\rho V_{g_2}=1/4`$ and $`\mathrm{\Pi }_{g_1}(\rho )=\rho `$. The corresponding roots are
$`(\underset{¯}{+,0^{n_01}},0^{n_0},\underset{¯}{+,0^{n_11}},0^{n_1})(0^{n_0},\underset{¯}{+,0^{n_01}},0^{n_1},\underset{¯}{+,0^{n_11}}),`$
$`(0^{n_0},\underset{¯}{,0^{n_01}},\underset{¯}{,0^{n_11}},0^{n_1})(\underset{¯}{,0^{n_01}},0^{n_0},0^{n_1},\underset{¯}{,0^{n_11}}),`$
$`(\underset{¯}{,0^{n_01}},0^{n_0},\underset{¯}{+,0^{n_11}},0^{n_1})(0^{n_0},\underset{¯}{,0^{n_01}},0^{n_1},\underset{¯}{+,0^{n_11}}),`$
$`(\underset{¯}{+,0^{n_01}},0^{n_0},0^{n_1},\underset{¯}{,0^{n_11}})(0^{n_0},\underset{¯}{+,0^{n_01}},\underset{¯}{,0^{n_11}},0^{n_1}).`$
They form a matter field transforming in the bifundamental $`(\text{ }\text{ }\text{ }\text{ }\text{ },\text{ }\text{ }\text{ }\text{ }\text{ })`$ representation of the gauge group. The matter fields from the third complex plane correspond to roots that satisfy $`\rho V_{g_2}=1/4`$ and $`\mathrm{\Pi }_{g_1}(\rho )=\rho `$. Again, one finds $`4n_0n_1`$ roots, giving a second bifundamental.
A similar analysis can be performed for the other boundary conditions, $`(+)`$, $`(++)`$, $`()`$. The result is shown in tables 36. One finds that the $`(+)`$ model and the $`()`$ model have non-Abelian gauge anomalies in the 33 sector. In the language of the previous subsection, these two boundary conditions correspond to $`\mu _{1}^{}{}_{}{}^{3}=1`$. There, we saw that precisely these models need $`D7`$-branes to cancel the tadpoles. A minimal choice consists in a set of $`D7_2`$ branes satisfying $`Tr\gamma _{(0,2),7_2}=16\mu _{3}^{}{}_{}{}^{3}\delta _{(0,1)}`$. Let us consider the $`()`$ model with $`\gamma _\mathrm{\Omega }`$ symmetric. According to (4.3), this gives $`\delta _g=1`$. From (c)), we have $`\mu _{1}^{}{}_{}{}^{7_2}=\mu _{1}^{}{}_{}{}^{3}`$ and $`\mu _{3}^{}{}_{}{}^{7_2}=\mu _{3}^{}{}_{}{}^{3}`$. Thus the theory on the $`D7_2`$-branes has $`(++)`$ boundary condition. As explained in , $`\gamma _{\mathrm{\Omega }_{7_2}}`$ has to be antisymmetric.<sup>10</sup><sup>10</sup>10If, instead of the standard $`\mathrm{\Omega }`$-projection discussed by Gimon and Polchinski , we use the alternative projection proposed by Dabholkar and Park in $`D=6`$ , $`\gamma _{\mathrm{\Omega },7_2}`$ and $`\gamma _{\mathrm{\Omega },3}`$ have the same symmetry. In table 3 the general solution for the gauge theory on a set of $`D`$-branes at a $`\text{/}\text{\_Z}_2\times \text{/}\text{\_Z}_4`$ singularity is shown for the case of symmetric $`\gamma _\mathrm{\Omega }`$. Knowing that changing $`\gamma _\mathrm{\Omega }`$ from symmetric to antisymmetric exchanges $`SO`$-factors with $`USp`$-factors, we find that the gauge theory on the $`D7_2`$-branes is $`USp(2m_0)\times USp(2m_1)`$. In order to cancel the tadpoles, we need $`2(2m_02m_1)=16`$. The $`37_2`$ sector gives matter fields transforming in the representations $`(\text{ }\text{ }\text{ }\text{ }\text{ };\text{ }\text{ }\text{ }\text{ }\text{ },1\mathrm{I})`$ and $`(\overline{\text{ }\text{ }\text{ }\text{ }\text{ }};1\mathrm{I},\text{ }\text{ }\text{ }\text{ }\text{ })`$ under the total gauge group $`U(n_0+n_1)\times USp(2m_0)\times USp(2m_1)`$. The total anomaly is thus given by $`2m_02m_18`$. We see that the condition of anomaly freedom is equivalent to the condition of vanishing tadpoles.
In the same way the $`\text{/}\text{\_Z}_2\times \text{/}\text{\_Z}_6`$ and the $`\text{/}\text{\_Z}_2\times \text{/}\text{\_Z}_6^{}`$ orientifold can be constructed. The result is shown in tables 36. Again, some of these models have non-Abelian gauge anomalies in the 33 sector. It can be seen that the same models require $`D7`$-branes for tadpole cancellation.
In general, one finds the following gauge groups for $`\text{/}\text{\_Z}_2\times \text{/}\text{\_Z}_N`$ orientifolds with discrete torsion (if $`\gamma _\mathrm{\Omega }`$ is symmetric):
$``$ boundary condition $`(++)`$:
$$\begin{array}{cc}N=4k+2:\hfill & G=SO(2n_0)\times _{i=1}^kU(n_i+n_{N/2i}),\hfill \\ N=4k:\hfill & G=SO(2n_0)\times SO(2n_k)\times _{i=1}^{k1}U(n_i+n_{N/2i}).\hfill \end{array}$$
(4.42)
$``$ boundary condition $`(+)`$:
$$\begin{array}{cc}N=4k+2:\hfill & G=SO(2n_k)\times _{i=0}^{k1}U(n_i+n_{N/21i}),\hfill \\ N=4k:\hfill & G=_{i=0}^{k1}U(n_i+n_{N/21i}).\hfill \end{array}$$
(4.43)
$``$ boundary condition $`(+)`$:
$$\begin{array}{cc}N=4k+2:\hfill & G=SO(2n_0)\times _{i=1}^kU(n_i+n_{N/2i}),\hfill \\ N=4k:\hfill & G=SO(2n_0)\times USp(2n_k)\times _{i=1}^{k1}U(n_i+n_{N/2i}).\hfill \end{array}$$
(4.44)
$``$ boundary condition $`()`$:
$$\begin{array}{cc}N=4k+2:\hfill & G=USp(2n_k)\times _{i=0}^{k1}U(n_i+n_{N/21i}),\hfill \\ N=4k:\hfill & G=_{i=0}^{k1}U(n_i+n_{N/21i}).\hfill \end{array}$$
(4.45)
### 4.4 Resolution of the singularities
Some of the deformations that are available in the orbifold case are absent in the orientifold case.
Take the easiest example: $`\text{/}\text{\_Z}_2\times \text{/}\text{\_Z}_2`$. In the orbifold case one can find three possible deformations :
$$\mathrm{\Delta }W=\underset{i=1}{\overset{3}{}}\zeta _iTrZ_i.$$
(4.46)
For the $`\text{/}\text{\_Z}_2\times \text{/}\text{\_Z}_2`$ with $`(++)`$ boundary condition, only one of the deformations survives the orientifold projection. Because of the symmetry of the $`Z_1`$, $`Z_3`$ (they are antisymmetric) and $`Z_2`$ (symmetric) matrices, one has:
$$\zeta _1TrZ_1=0,\zeta _3TrZ_3=0.$$
(4.47)
This indicates that the orientifold can only be deformed to the $`\text{/}\text{\_Z}_2`$ singularity. The additional deformation leading to the conifold is frozen in the orientifold case. Something similar happened in the $`\text{/}\text{\_Z}_2\times \text{/}\text{\_Z}_2`$ orientifold without discrete torsion .
## 5 Conclusions
We have seen that only real values of the discrete torsion parameter $`ϵ`$ are allowed for orientifold models, in contrast to the orbifold case, where $`ϵ`$ can take complex values. This condition can be understood from several viewpoints. In the open string sector this is related to the fact that the $`\gamma `$-matrices form real projective representations. These are classified by $`H^2(\mathrm{\Gamma },\mathrm{IR}^{})=H^2(\mathrm{\Gamma },\text{/}\text{\_Z}_2)`$, in a similar manner as complex projective representations are characterized by $`H^2(\mathrm{\Gamma },\mathrm{C}\text{ }^{})=H^2(\mathrm{\Gamma },U(1))`$. As a consequence, only $`ϵ=\pm 1`$ is allowed for orientifolds. In the closed string spectrum, the condition of real $`ϵ`$ is related to the matching between left and right moving degrees of freedom. In general, this matching is impossible if the Hodge numbers $`h^{1,2}`$ and $`h^{2,1}`$ from one twisted sector are different. One finds that $`h^{1,2}=h^{2,1}`$ is only guaranteed in each twisted sector if $`ϵ=\pm 1`$. Finally, one finds an inconsistency in the tadpole cancellation conditions for non-real $`ϵ`$. The characters of the projective representation $`\gamma `$ must have a precise value to cancel the Klein bottle contribution. This conditions cannot be satisfied if arbitrary values of the discrete torsion are allowed.
For the $`\text{/}\text{\_Z}_N\times \text{/}\text{\_Z}_M`$ orientifold non-trivial discrete torsion is only possible if $`N`$ and $`M`$ are both even. In this case, there are four non-equivalent orientifold models. This classification is based on the possibility of imposing different boundary conditions (‘vector structures’) on the $`\gamma `$-matrices. The tadpole conditions are different for each case. Some of them require $`D7`$-branes for consistency. This condition is equivalent to the requirement that non-Abelian gauge anomalies be absent.
We have also analyzed the resolution of the $`\text{/}\text{\_Z}_2\times \text{/}\text{\_Z}_2`$ singularity in both cases, orbifold and orientifold. In the orbifold case, one can get a $`\text{/}\text{\_Z}_2`$ singularity and a conifold by the deformations of the F-flatness conditions. In the orientifold only the $`\text{/}\text{\_Z}_2`$ singularity can be obtained.
This discussion can be generalized to compact orientifolds with discrete torsion . The S-dual heterotic models are related to higher level Kac-Moody algebras. In particular, as the only non-trivial discrete torsion is $`ϵ=1`$, we will find heterotic duals realized at Kac-Moody level 2 .
Acknowledgements
It is a pleasure to thank Angel Uranga and Luis Ibáñez for many helpful ideas and comments on the manuscript. The work of M.K. is supported by a TMR network of the European Union, ref. FMRX-CT96-0090. The work of R.R. is supported by the MEC through a FPU Grant.
Appendix
## Appendix A Projective representations
The matrices that represent the action of the orbifold group on the Chan-Paton indices of open strings form projective representations. In this appendix, we review some useful facts about the latter (for an introduction to this subject see ). Let $`\mathrm{\Gamma }`$ be a finite group. A projective representation is a mapping $`\mathrm{\Gamma }GL(n,\mathrm{IK})`$, which associates to each element $`g\mathrm{\Gamma }`$ a matrix $`\gamma _gGL(n,\mathrm{IK})`$ that satisfies the conditions<sup>11</sup><sup>11</sup>11Equivalently a projective representations can be defined as a homomorphism $`\mathrm{\Gamma }PGL(n,\mathrm{IK})`$.
$$\gamma _e=1\mathrm{I},\gamma _g\gamma _h=\alpha _{g,h}\gamma _{gh}g,h\mathrm{\Gamma },$$
(A.1)
where $`e`$ is the neutral element of $`\mathrm{\Gamma }`$ and $`\alpha _{g,h}\mathrm{IK}^{}`$ are arbitrary non-zero numbers. These numbers are called the factor system of the projective representation $`\gamma `$. Using associativity, one immediately obtains
$$\alpha _{g,e}=\alpha _{e,g}=1,\alpha _{g,hk}\alpha _{h,k}=\alpha _{g,h}\alpha _{gh,k}g,h,k\mathrm{\Gamma }.$$
(A.2)
A map $`\alpha :\mathrm{\Gamma }\times \mathrm{\Gamma }\mathrm{IK}^{}`$, $`(g,h)\alpha _{g,h}`$, with the above properties is called a cocycle. The set of all cocycles is denoted by $`Z^2(\mathrm{\Gamma },\mathrm{IK}^{})`$. A projective representation $`\widehat{\gamma }`$ is considered equivalent to $`\gamma `$ if it is obtained by substituting $`\widehat{\gamma }_g=\rho _g\gamma _g`$, with $`\rho _g\mathrm{IK}^{}`$. The cocycles are then related by
$$\widehat{\alpha }_{g,h}=\alpha _{g,h}\rho _g\rho _h\rho _{gh}^1.$$
(A.3)
This motivates the following definition. Given a map $`\rho :\mathrm{\Gamma }\mathrm{IK}^{}`$, $`g\rho _g`$, with $`\rho _e=1`$, one defines the associated coboundary by $`\delta \rho :\mathrm{\Gamma }\times \mathrm{\Gamma }\mathrm{IK}^{}`$, $`(g,h)\rho _g\rho _h\rho _{gh}^1`$. The set of all coboundaries, $`B^2(\mathrm{\Gamma },\mathrm{IK}^{})`$, is a subset of $`Z^2(\mathrm{\Gamma },\mathrm{IK}^{})`$. Two cocycles are equivalent if they differ only by a coboundary. We see that the non-equivalent projective representations are characterized by the elements of $`H^2(\mathrm{\Gamma },\mathrm{IK}^{})=Z^2(\mathrm{\Gamma },\mathrm{IK}^{})/B^2(\mathrm{\Gamma },\mathrm{IK}^{})`$.
If $`\mathrm{\Gamma }`$ is Abelian, one finds from (A.1) that
$$\gamma _g\gamma _h=\beta _{g,h}\gamma _h\gamma _g,\mathrm{where}\beta _{g,h}=\alpha _{g,h}\alpha _{h,g}^1.$$
(A.4)
The $`\beta _{g,h}`$ only depend on the equivalence class of the $`\alpha _{g,h}`$. Furthermore they satisfy
$$\beta _{g,g}=1,\beta _{g,h}=\beta _{h,g}^1,\beta _{g,hk}=\beta _{g,h}\beta _{g,k}g,h,k\mathrm{\Gamma }.$$
(A.5)
It is clear that to each element $`[\alpha ]`$ of $`H^2(\mathrm{\Gamma },\mathrm{IK}^{})`$ there corresponds a unique cocycle<sup>12</sup><sup>12</sup>12Note that $`\beta `$ is a cocycle because it satisfies (A.2), but it is not a factor system because its definition differs from (A.1). $`\beta `$, as given in (A.4). On the other hand to each $`\beta `$ there corresponds a unique (up to equivalence) factor system $`\alpha _{g,h}=\sqrt{\beta _{g,h}}`$. As a consequence, the projective representations of an Abelian finite group $`\mathrm{\Gamma }`$ can be characterized either by (A.1), (A.2) or by the first eq. of (A.4) and the three eqs. of (A.5). Both descriptions are equivalent up to a transformation (A.3). Sometimes it is useful to fix as many of the $`\alpha _{g,h}`$ as possible without putting any restriction on the equivalence class $`[\alpha ]`$. A convenient choice is:
$$\alpha _{g_1^a,g_2^b}=\alpha _{g_1,g_1^a}=\alpha _{g_2,g_2^b}=1,a=1,\mathrm{},N,b=1,\mathrm{},M.$$
(A.6)
This corresponds to choosing a set of matrices that satisfies
$$(\gamma _{g_1})^a(\gamma _{g_2})^b=\gamma _{g_1^ag_2^b}.$$
(A.7)
For physical reasons we restrict ourselves to unitary projective representations over the complex or real numbers, i.e. $`\gamma _g\gamma _g^{}=1\mathrm{I}`$ and $`\mathrm{IK}=\mathrm{C}\text{ }`$ or $`\mathrm{IK}=\mathrm{IR}`$. One can prove that<sup>13</sup><sup>13</sup>13For the cocycles related to unitary projective representations this is obvious. In general, this is a consequence of proposition 2.3.10 and lemma 2.3.19 of .
$$H^2(\mathrm{\Gamma },\mathrm{C}\text{ }^{})=H^2(\mathrm{\Gamma },U(1)),H^2(\mathrm{\Gamma },\mathrm{IR}^{})=H^2(\mathrm{\Gamma },\text{/}\text{\_Z}_2).$$
(A.8)
In this article we will be mainly interested in the case where $`\mathrm{\Gamma }=\text{/}\text{\_Z}_N\times \text{/}\text{\_Z}_M`$. Let $`g_1`$ and $`g_2`$ be the generators of $`\text{/}\text{\_Z}_N`$ and $`\text{/}\text{\_Z}_M`$ respectively and denote the element $`g_1^ag_2^b\mathrm{\Gamma }`$ by $`(a,b)`$. The form of the $`\beta `$ cocycles is completely fixed by the conditions (A.5). If $`\mathrm{IK}=\mathrm{C}\text{ }`$, there are $`p=\mathrm{gcd}(N,M)`$ non-equivalent cocycles, given by
$$\beta _{(a,b),(a^{},b^{})}=\omega _p^{m(ab^{}ba^{})},\mathrm{with}\omega _p=e^{2\pi i/p},m=1,\mathrm{},p.$$
(A.9)
The different projective representations are therefore determined by the parameter $`ϵ=\omega _p^m`$. Choosing the factor system (A.6), one finds a simple relation for the product of two $`\gamma `$-matrices:
$$\gamma _{(a,b)}\gamma _{(c,d)}=ϵ^{bc}\gamma _{(a+c,b+d)}.$$
(A.10)
The possibility of having $`ϵ1`$ is the open string analogue of the discrete torsion in closed string orbifolds discussed in . The open string models with discrete torsion are distinguished by the integer $`s=p/\mathrm{gcd}(m,p)`$ (this is the smallest non-zero number such that $`ϵ^s=1`$). According to , any irreducible projective representation $`_{\mathrm{irr}}^{(s)}`$ is $`s`$-dimensional and (up to projective equivalence) of the form
$$\gamma _{g_1}=\mathrm{diag}(1,ϵ^1,ϵ^2,\mathrm{},ϵ^{(s1)}),\gamma _{g_2}=\left(\begin{array}{ccccc}0& 1& 0& \mathrm{}& 0\\ 0& 0& 1& \mathrm{}& \\ \mathrm{}& & \mathrm{}& \mathrm{}& \\ 0& & \mathrm{}& 0& 1\\ 1& 0& \mathrm{}& & 0\end{array}\right).$$
(A.11)
All the the irreducible projective representations $`_{\mathrm{irr},k,l}^{(s)}`$ that are linearly non-equivalent (i.e. they belong to different factor systems but, of course, are projectively equivalent) can be obtained by multiplying the matrices in (A.11) by phases:
$$\widehat{\gamma }_{g_1}=\omega _N^k\gamma _{g_1},k=0,\mathrm{},\frac{N}{s}1,\widehat{\gamma }_{g_2}=\omega _M^l\gamma _{g_1},l=0,\mathrm{},\frac{M}{s}1.$$
(A.12)
A general projective representation $`^{(s)}`$ is a direct sum of irreducible blocks $`_{\mathrm{irr},k,l}^{(s)}`$:
$$^{(s)}=\underset{k,l}{}n_{kl}_{\mathrm{irr},k,l}^{(s)}.$$
(A.13)
This representation has dimension $`s_{k,l}n_{kl}`$. The regular represention, of dimension $`|\mathrm{\Gamma }|=NM`$, is obtained by setting $`n_{kl}=sk,l`$. Let $`\gamma _g^{(s)}`$ denote the matrix associated to the group element $`g`$ in the general projective representation $`^{(s)}`$. The character of $`g`$ in $`^{(s)}`$ is defined to be the trace of this matrix. Denoting again $`g=g_1^ag_2^b`$ by $`(a,b)`$ we find
$$\chi ^{(^{(s)})}(g)Tr\gamma _g^{(s)}=\{\begin{array}{cc}0\hfill & \mathrm{if}(a,b)s\text{/}\text{\_Z}\times s\text{/}\text{\_Z}\hfill \\ s_{k,l}n_{k,l}\omega _N^{ka}\omega _M^{lb}\hfill & \mathrm{if}(a,b)s\text{/}\text{\_Z}\times s\text{/}\text{\_Z}\hfill \end{array}$$
(A.14)
The solutions for $`\mathrm{IK}=\mathrm{IR}`$ can be obtained by restricting $`\beta `$ to be real. The non-equivalent cocycles correspond to setting $`m=p/2`$ (if $`p`$ is even) or $`m=p`$ in (A.9). We see that discrete torsion is only possible if $`p`$ is even. The only non-trivial cocycle in this case is $`\beta _{(a,b),(a^{},b^{})}=(1)^{ab^{}ba^{}}`$. The irreducible projective representations are again of the form (A.11), with $`ϵ=1`$ and $`s=2`$. However, there are four different complex projective representations which are not equivalent over the real numbers and which have a real factor system<sup>14</sup><sup>14</sup>14We used that $`M`$ and $`N`$ are even, which is true because $`\mathrm{gcd}(M,N)`$ is even.:
$`(\mathrm{i})\gamma _{g_1},\gamma _{g_2},`$
$`(\mathrm{ii})\omega _{2N}\gamma _{g_1},\gamma _{g_2},`$ (A.15)
$`(\mathrm{iii})\gamma _{g_1},\omega _{2M}\gamma _{g_2},`$
$`(\mathrm{iv})\omega _{2N}\gamma _{g_1},\omega _{2M}\gamma _{g_2},`$
where $`\gamma _{g_{1/2}}`$ are given in (A.11). In general, an irreducible complex projective representation is of the form $`\stackrel{~}{\gamma }_{g_{1/2}}=\rho _{1/2}\gamma _{g_{1/2}}`$, with $`\rho _{1/2}\mathrm{C}\text{ }^{}`$. The condition that the factor system be real implies $`(\stackrel{~}{\gamma }_{g_1})^N=c_1\mathrm{\hspace{0.17em}1}\mathrm{I}`$, $`(\stackrel{~}{\gamma }_{g_2})^M=c_2\mathrm{\hspace{0.17em}1}\mathrm{I}`$, with $`c_{1/2}\mathrm{IR}^{}`$. Therefore, $`\rho _1^N=\pm 1`$ and $`\rho _2^M=\pm 1`$. These are the four possibilities of (A). Defining $`\eta =(\eta _1,\eta _2)`$, with $`\eta _i\{0,1\}`$ by
$$(\gamma _{g_1})^N=(1)^{\eta _1}1\mathrm{I},(\gamma _{g_2})^M=(1)^{\eta _2}1\mathrm{I},$$
(A.16)
we can give the form of the linearly non-equivalent real projective representations:
$$_{\mathrm{irr},k,l,\eta }^{\mathrm{real}}=\{\begin{array}{cc}_{\mathrm{irr},k,l,\eta }^{(2)}\hfill & \mathrm{if}k=l=\eta _1=\eta _2=0\hfill \\ _{\mathrm{irr},k,l,\eta }^{(2)}(_{\mathrm{irr},k,l,\eta }^{(2)})^{}\hfill & \mathrm{else}\hfill \end{array},$$
(A.17)
where $`_{\mathrm{irr},k,l}^{(2)}`$ is obtained from (A.11) by multiplying $`\gamma _{g_1}`$ with phases $`\omega _N^k\omega _{2N}^{\eta _1}`$ and $`\gamma _{g_2}`$ with phases $`\omega _M^l\omega _{2M}^{\eta _2}`$ as in (A.12) and (A). These projective representations are either two-dimensional (first line of (A.17)) or four-dimensional (second line of (A.17)). Again, a general projective representation $`_\eta ^{\mathrm{real}}`$ is a direct sum of irreducible blocks $`_{\mathrm{irr},k,l,\eta }^{\mathrm{real}}`$:
$$_\eta ^{\mathrm{real}}=\underset{k,l}{}n_{kl}_{\mathrm{irr},k,l,\eta }^{\mathrm{real}}.$$
(A.18)
This representation has dimension $`2rn_{0,0}+4_{(k,l)(0,0)}n_{kl}`$, where $`r=2^{|\eta _1\eta _2+\eta _1\eta _2|}`$. The regular representation is obtained by setting $`n_{0,0}=2/r`$ and $`n_{kl}=1(k,l)(0,0)`$. For the character of an element $`g=g_1^ag_2^b`$ we find
$`\chi ^{(_\eta ^{\mathrm{real}})}(g)`$ $``$ $`Tr\gamma _g^{\mathrm{real}}`$ (A.19)
$`=`$ $`\{\begin{array}{cc}0\hfill & \text{if }a\text{ or }b\text{ is odd}\hfill \\ 2rn_{0,0}+4_{(k,l)(0,0)}n_{k,l}\mathrm{Re}(\omega _N^{(k+\eta _1/2)a}\omega _M^{(l+\eta _2/2)b})\hfill & \text{if }a\text{ and }b\text{ are even}\hfill \\ \text{with }r=1\text{ if }\eta _1=\eta _2=0\text{ and }r=2\text{ else}\hfill & \end{array}`$ (A.23)
## Appendix B Tadpoles for $`𝐙_N\times 𝐙_M`$ orientifolds
We sketch the calculation of the tadpoles for non-compact $`\text{/}\text{\_Z}_N\times \text{/}\text{\_Z}_M`$ orientifolds. Basically, we follow the work of .
The set-up consists of several $`D3`$-branes at the orbifold singularity plus some $`D7_i`$-branes, where the $`i`$ denotes the complex plane with Dirichlet conditions. We take the orientifold involution $`\mathrm{\Omega }^{}=\mathrm{\Omega }(1)^{F_L}R_1R_2R_3`$ (other types can be considered, as in ).
The procedure we will follow is based on the computation of the tadpoles for $`T^6/(\text{/}\text{\_Z}_N\times \text{/}\text{\_Z}_M)`$ and then taking the non-compact limit. In order to take this limit, one must ignore the untwisted tadpoles and the contributions with an inversely proportional dependence on the volume. The dependence proportional to the volume is related to some factors from the momentum modes which become continuous in the non-compact limit.
To simplify the notation, let us define $`s_i=\mathrm{sin}(\pi \overline{k}\overline{v}_i)`$, $`c_i=\mathrm{cos}(\pi \overline{k}\overline{v}_i)`$ and $`\stackrel{~}{s}_i=\mathrm{sin}(2\pi \overline{k}\overline{v}_i)`$. We will not give explicitly the volume dependence, but this dependence can be extracted from the zeros and divergences of each contribution.
The cylinder amplitude can be split into four sectors: $`33`$, $`7_i7_i`$, $`7_i7_j`$ and $`7_i3`$. All the contributions can be neatly recast:
$$𝒞=\underset{\overline{k}=(1,1)}{\overset{(N,M)}{}}\frac{1}{8s_1s_2s_3}\left[8s_1s_2s_3Tr\gamma _{\overline{k},3}+_{i=1}^32s_iTr\gamma _{\overline{k},7_i}\right]^2.$$
(B.1)
Note that in the orbifold case, this is the only contribution, that has to be considered.
Because of the $`J`$ operation that relates the $`\overline{k}`$ with the $`\overline{k}`$ twisted sector, the contributions to the Klein bottle come from the untwisted and order-two twisted sectors. Three cases must be distinguished:
* $`N`$ and $`M`$ are both odd: There are no order-two twisted sectors.
* $`N`$ is odd and $`M`$ is even: There is only one order-two sector: $`\overline{k}_1=(0,M/2)`$. It fixes the first complex plane. Note that in this case no discrete torsion is allowed. The Klein bottle contribution is of the form:
$$𝒦=\underset{\overline{k}}{}\left(𝒦_0(\overline{k})+𝒦_1(\overline{k})\right),$$
(B.2)
where $`𝒦_0`$ is the contribution from the untwisted sector and $`𝒦_1`$ is the contribution from the order-two twisted sector that fixes the first complex plane.
* $`N`$ and $`M`$ are both even: There are three order-two sectors: $`\overline{k}_3=(N/2,0)`$, $`\overline{k}_1=(0,M/2)`$ and $`\overline{k}_2=(N/2,M/2)`$. The Klein bottle contribution can be written as
$$𝒦=\underset{\overline{k}}{}\left(𝒦_0(\overline{k})+𝒦_1(\overline{k})+𝒦_2(\overline{k})+𝒦_3(\overline{k})\right),$$
(B.3)
where $`𝒦_i`$ is the contribution from the order-two twist that fixes the $`i`$-th complex plane. The untwisted sector contribution is of the form
$$𝒦_0(\overline{k})=16\frac{2\stackrel{~}{s}_12\stackrel{~}{s}_22\stackrel{~}{s}_3}{4c_1^24c_2^24c_3^2},$$
(B.4)
and the order-two twisted contributions are
$$𝒦_i(\overline{k})=16\stackrel{~}{ϵ}_i\beta _i\frac{2\stackrel{~}{s}_i}{4c_i^2},$$
(B.5)
where $`\stackrel{~}{ϵ}_i`$ is a sign that weights the contribution of the sector that fixes the $`i`$-th complex plane and $`\beta _i=\beta _{\overline{k}_i,\overline{k}}`$.
This Klein bottle contributions can be rewritten for each of the three cases:
* $`N`$ odd, $`M`$ odd:
$$𝒦=\underset{\overline{k}=(1,1)}{\overset{(N,M)}{}}\frac{1}{8\stackrel{~}{s}_1\stackrel{~}{s}_2\stackrel{~}{s}_3}[32s_1s_2s_3]^2.$$
(B.6)
* $`N`$ odd, $`M`$ even:
$$𝒦=\underset{\overline{k}=(1,1)}{\overset{(N,M/2)}{}}\frac{1}{8\stackrel{~}{s}_1\stackrel{~}{s}_2\stackrel{~}{s}_3}[32(s_1s_2s_3\stackrel{~}{ϵ}_1s_1c_2c_3)]^2.$$
(B.7)
* $`N`$ even, $`M`$ even:
$$𝒦=\underset{\overline{k}=(1,1)}{\overset{(N/2,M/2)}{}}\frac{1}{8\stackrel{~}{s}_1\stackrel{~}{s}_2\stackrel{~}{s}_3}[32(s_1s_2s_3\underset{ijk}{}\stackrel{~}{ϵ}_i\beta _is_ic_jc_k)]^2.$$
(B.8)
Finally, one must consider the Möbius strip contribution. For each of the four types of branes there is a contribution:
$$_3(\overline{k})=88s_1s_2s_3Tr\left(\gamma _{\mathrm{\Omega }\overline{k},3}^1\gamma _{\mathrm{\Omega }\overline{k},3}^{}\right),$$
(B.9)
$$_{7_i}(\overline{k})=8\frac{8s_ic_jc_k}{4c_j^24c_k^2}Tr\left(\gamma _{\mathrm{\Omega }\overline{k},7_i}^1\gamma _{\mathrm{\Omega }\overline{k},7_i}^{}\right).$$
(B.10)
Using the properties (4.21), (4.22), (A.10) of projective representations, we find
$$\gamma _{\mathrm{\Omega }\overline{k}}^1\gamma _{\mathrm{\Omega }\overline{k}}^{}=\delta _{\overline{k}}(\gamma _{\overline{k}}^2)^{}=ϵ^{k_1k_2}\delta _{\overline{k}}(\gamma _{\overline{2}k})^{},$$
(B.11)
where $`\overline{k}=(k_1,k_2)`$ and we chose a factor system such that $`(\gamma _{(1,0)})^{k_1}(\gamma _{(0,1)})^{k_2}=\gamma _{\overline{k}}`$.
Reordering the contributions, as we have done for the Klein bottle, one can obtain:
* $`N`$ odd, $`M`$ odd:
From $`D3`$-branes:
$$_3=2\underset{\overline{k}=(1,1)}{\overset{(N,M)}{}}\frac{\delta _{\overline{k},3}}{8\stackrel{~}{s}_1\stackrel{~}{s}_2\stackrel{~}{s}_3}[32s_1s_2s_3][8\stackrel{~}{s}_1\stackrel{~}{s}_2\stackrel{~}{s}_3Tr\gamma _{2\overline{k},3}].$$
(B.12)
From $`D7_i`$-branes:
$$_{7_i}=2\underset{\overline{k}=(1,1)}{\overset{(N,M)}{}}\frac{\delta _{\overline{k},7_i}}{8\stackrel{~}{s}_1\stackrel{~}{s}_2\stackrel{~}{s}_3}[32s_1s_2s_3][2\stackrel{~}{s}_iTr\gamma _{2\overline{k},7_i}].$$
(B.13)
* $`N`$ odd, $`M`$ even:
From $`D3`$-branes:
$$_3=2\underset{\overline{k}=(1,1)}{\overset{(N,M/2)}{}}\frac{\delta _{\overline{k},3}}{8\stackrel{~}{s}_1\stackrel{~}{s}_2\stackrel{~}{s}_3}[32(s_1s_2s_3\stackrel{~}{ϵ}_1s_1c_2c_3)][8\stackrel{~}{s}_1\stackrel{~}{s}_2\stackrel{~}{s}_3Tr\gamma _{2\overline{k},3}].$$
(B.14)
From $`D7_i`$-branes:
$$_{7_i}=2\underset{\overline{k}=(1,1)}{\overset{(N,M/2)}{}}\frac{\delta _{\overline{k},7_i}}{8\stackrel{~}{s}_1\stackrel{~}{s}_2\stackrel{~}{s}_3}[32(s_1s_2s_3\stackrel{~}{ϵ}_1s_1c_2c_3)][2\stackrel{~}{s}_iTr\gamma _{2\overline{k},7_i}].$$
(B.15)
This factorization is necessary to match the cylinder and the Klein bottle contributions. But to obtain it, we had to impose some restrictions on the ‘vector structures’ $`\mu _i`$ (introduced in (4.3)) and the signs $`\stackrel{~}{ϵ}_i`$:
$$\mu _{1}^{}{}_{}{}^{3}=\mu _{1}^{}{}_{}{}^{7_1}=\mu _{1}^{}{}_{}{}^{7_2}=\mu _{1}^{}{}_{}{}^{7_3}=\stackrel{~}{ϵ}_1.$$
(B.16)
Here we took $`\delta _{\overline{k}}=1\overline{k}`$ if $`\gamma _\mathrm{\Omega }`$ is symmetric and $`\delta _{\overline{k}}=1\overline{k}`$ if $`\gamma _\mathrm{\Omega }`$ is antisymmetric, as in (4.3).
* $`N`$ even, $`M`$ even:
From $`D3`$-branes:
$$_3=2\underset{\overline{k}=(1,1)}{\overset{(N/2,M/2)}{}}\frac{ϵ^{k_1k_2}\delta _{\overline{k},3}}{8\stackrel{~}{s}_1\stackrel{~}{s}_2\stackrel{~}{s}_3}[32(s_1s_2s_3\underset{ijk}{}\stackrel{~}{ϵ}_i\beta _is_ic_jc_k)][8\stackrel{~}{s}_1\stackrel{~}{s}_2\stackrel{~}{s}_3Tr\gamma _{2\overline{k},3}].$$
(B.17)
From $`D7_i`$-branes:
$$_{7_i}=2\underset{\overline{k}=(1,1)}{\overset{(N/2,M/2)}{}}\frac{ϵ^{k_1k_2}\delta _{\overline{k},7_i}}{8\stackrel{~}{s}_1\stackrel{~}{s}_2\stackrel{~}{s}_3}[32(s_1s_2s_3\underset{ijk}{}\stackrel{~}{ϵ}_i\beta _is_ic_jc_k)][2\stackrel{~}{s}_iTr\gamma _{2\overline{k},7_i}].$$
(B.18)
Again, this factorization is necessary to match the cylinder and the Klein bottle contributions, but it is only possible if we impose some restrictions on the ‘vector structures’ $`\mu _i`$ and the signs $`\stackrel{~}{ϵ}_i`$:
$`\mu _{1}^{}{}_{}{}^{3}=\mu _{1}^{}{}_{}{}^{7_1}=\mu _{1}^{}{}_{}{}^{7_2}=\mu _{1}^{}{}_{}{}^{7_3}`$ $`=`$ $`\stackrel{~}{ϵ}_1,`$
$`\mu _{3}^{}{}_{}{}^{3}=\mu _{3}^{}{}_{}{}^{7_1}=\mu _{3}^{}{}_{}{}^{7_2}=\mu _{3}^{}{}_{}{}^{7_3}`$ $`=`$ $`\stackrel{~}{ϵ}_3,`$ (B.19)
$`\stackrel{~}{ϵ}_1\stackrel{~}{ϵ}_2\stackrel{~}{ϵ}_3`$ $`=`$ $`ϵ^{MN/4}.`$
One further condition is needed, to be able to write the total tadpole contribution $`𝒞++𝒦`$ as a square<sup>15</sup><sup>15</sup>15This condition is necessary if one uses the standard $`\mathrm{\Omega }`$-projection of . Using, however, the alternative projection of , this condition can be relaxed.:
$$\delta _{\overline{k},3}=\delta _{\overline{k},7_1}=\delta _{\overline{k},7_2}=\delta _{\overline{k},7_3}\delta _{\overline{k}}.$$
(B.20)
If the above conditions are satisfied, the tadpole conditions can be written in the form that appears in the main text, eqs. (4.7)–(4.12).
## Appendix C Shift formalism and discrete torsion
Let us sketch a general procedure for obtaining the open string spectrum of an arbitrary type IIB orbifold or orientifold with discrete torsion. The method is based on the relation between the shift formalism and the one using matrices.
The open string spectrum of a general type IIB orbifold with or without discrete torsion is determined by the matrices $`\lambda `$ that satisfy
$$\gamma _g\lambda \gamma _g^1=r(g)\lambda ,$$
(C.1)
where $`\gamma _g`$ represents the action of the orbifold group $`\mathrm{\Gamma }`$ on the Chan-Paton indices of the open string and $`r(g)`$ the action of $`\mathrm{\Gamma }`$ on the oscillator state of the open string. For the gauge degrees of freedom, $`\lambda ^{(0)}`$, one has $`r(g)=1`$, whereas $`r(g)=e^{2\pi iv_i}`$ for the matter degrees of freedom associated to the $`i`$-th complex plane, $`\lambda ^{(i)}`$, if $`g`$ acts as $`z_ie^{2\pi iv_i}z_i`$ on the $`i`$-th complex coordinate. A straightforward but cumbersome method to get the the spectrum, consists in constructing explicitely the matrices $`\lambda ^{(0)}`$, $`\lambda ^{(i)}`$ that solve (C.1). The $`\lambda ^{(0)}`$ transform in the adjoint representation of the gauge group $`G`$. Knowing $`G`$, one finds the matter representations by looking at the transformation properties of the $`\lambda ^{(i)}`$.
In it has been shown that this is equivalent to a shift formalism which is very similar to the one that is known from heterotic orbifolds. If all the matrices $`\gamma _g`$ commute, one can diagonalize them simultaneously by choosing the Cartan-Weyl basis:
$$\gamma _g=e^{2\pi iV_gH},$$
(C.2)
where $`H=(H_1,H_2,\mathrm{},H_{rank(\stackrel{~}{G})})`$ is a vector containing the Cartan generators of the gauge group $`\stackrel{~}{G}`$ (before the orbifold projection) and $`V_g=(V_1,V_2,\mathrm{},V_{rank(\stackrel{~}{G})})`$ is the shift vector that represents the action of $`g\mathrm{\Gamma }`$ on the root lattice of $`\stackrel{~}{G}`$. One can show that in the Cartan-Weyl basis the identity
$$\gamma _gE_a\gamma _g^1=e^{2\pi i\rho ^aV_g}E_a,a=1,\mathrm{},dim(\stackrel{~}{G})rank(\stackrel{~}{G}),$$
(C.3)
holds for all charged generators $`E_a`$ of $`\stackrel{~}{G}`$, if $`\rho ^a`$ is the root vector associated to $`E_a`$, i.e. $`[H_I,E_a]=\rho _I^aE_a`$. The condition (C.1) now reads
$`\rho ^aV_g`$ $`=`$ $`0\mathrm{mod}\text{/}\text{\_Z}\text{for gauge fields},`$ (C.4)
$`\rho ^aV_g`$ $`=`$ $`v_i\mathrm{mod}\text{/}\text{\_Z}\text{for matter fields from the }i\text{-th complex plane}.`$ (C.5)
In this representation, it is very easy to find the gauge group $`G\stackrel{~}{G}`$ that is preserved after the orbifold projection and the representations of the matter fields.
If $`\mathrm{\Gamma }=\text{/}\text{\_Z}_N\times \text{/}\text{\_Z}_M`$ and the projective representation $`\gamma `$ has discrete torsion, then $`\gamma _{g_1}\gamma _{g_2}\gamma _{g_2}\gamma _{g_1}`$, where $`g_1,g_2`$ are the generators of $`\text{/}\text{\_Z}_N,\text{/}\text{\_Z}_M`$. The above method must be modified because $`\gamma _{g_1}`$ and $`\gamma _{g_2}`$ cannot be simultaneously diagonalized. Following the idea of , we diagonalize $`\gamma _{g_1}`$ and represent $`g_2`$ by a permutation acting on the entries of the root vectors.
For a $`\text{/}\text{\_Z}_N\times \text{/}\text{\_Z}_M`$ orbifold with discrete torsion $`ϵ`$, the matrices $`\gamma _{g_1}`$, $`\gamma _{g_2}`$ are of the following form (see appendix A):
$`\gamma _{g_1}`$ $`=`$ $`{\displaystyle \underset{k}{}}(\omega _N^k\gamma _{g_1}^{\mathrm{irrep}})1\mathrm{I}_{n_k},\gamma _{g_2}={\displaystyle \underset{l}{}}(\omega _M^l\gamma _{g_2}^{\mathrm{irrep}})1\mathrm{I}_{n_l},`$ (C.6)
$`\mathrm{with}`$ $`\gamma _{g_1}^{\mathrm{irrep}}=\mathrm{diag}(1,ϵ^1,ϵ^2,\mathrm{},ϵ^{(s1)}),\gamma _{g_2}^{\mathrm{irrep}}=\left(\begin{array}{ccccc}0& 1& 0& \mathrm{}& 0\\ 0& 0& 1& \mathrm{}& \\ \mathrm{}& & \mathrm{}& \mathrm{}& \\ 0& & \mathrm{}& 0& 1\\ 1& 0& \mathrm{}& & 0\end{array}\right),`$
$`\omega _p=e^{2\pi i/p},k=0,\mathrm{},{\displaystyle \frac{N}{s}}1,l=0,\mathrm{},{\displaystyle \frac{M}{s}}1,`$
$`s\text{is the minimal positive integer, such that }ϵ^s=1.`$
In the Cartan-Weyl basis, $`\gamma _{g_1}`$ corresponds to the shift
$`V_{g_1}`$ $`=`$ $`{\displaystyle \frac{1}{N}}(0^{n_0},\left({\displaystyle \frac{N}{s}}\right)^{n_0},\mathrm{},\left((s1){\displaystyle \frac{N}{s}}\right)^{n_0},1^{n_1},({\displaystyle \frac{N}{s}}+1)^{n_1},\mathrm{},((s1){\displaystyle \frac{N}{s}}+1)^{n_1},`$ (C.13)
$`\mathrm{},({\displaystyle \frac{N}{s}}1)^{n_{N/s1}},(2{\displaystyle \frac{N}{s}}1)^{n_{N/s1}},\mathrm{},(N1)^{n_{N/s1}}).`$
The matrix $`\gamma _{g_2}`$ acts as a permutation $`\mathrm{\Pi }_{g_2}`$ on the root vectors $`\rho `$. We assume $`NM`$ and consider first the case $`s=M`$, i.e. $`\gamma _{g_2}=\gamma _{g_2}^{\mathrm{irrep}}1\mathrm{I}_n`$, where $`n=_kn_k`$:
$`\mathrm{\Pi }_{g_2}:`$ $`(\stackrel{~}{\rho }_1^{(n_0)},\mathrm{},\stackrel{~}{\rho }_s^{(n_0)},\stackrel{~}{\rho }_1^{(n_1)},\mathrm{},\stackrel{~}{\rho }_s^{(n_1)},\mathrm{},\stackrel{~}{\rho }_1^{(n_{N/s1})},\mathrm{},\stackrel{~}{\rho }_s^{(n_{N/s1})})`$
$`(\stackrel{~}{\rho }_s^{(n_0)},\stackrel{~}{\rho }_1^{(n_0)},\mathrm{},\stackrel{~}{\rho }_{s1}^{(n_0)},\stackrel{~}{\rho }_s^{(n_1)},\stackrel{~}{\rho }_1^{(n_1)},\mathrm{},\stackrel{~}{\rho }_{s1}^{(n_1)},\mathrm{},\stackrel{~}{\rho }_s^{(n_{N/s1})},\mathrm{},\stackrel{~}{\rho }_{s1}^{(n_{N/s1})}),`$
$`\mathrm{where}`$ $`\stackrel{~}{\rho }_i^{(n_k)}=(\rho _{s(n_0+\mathrm{}+n_{k1})+(i1)n_k+1},\mathrm{},\rho _{s(n_0+\mathrm{}+n_{k1})+in_k}).`$
The gauge group $`G`$ of the orbifold model is determined by finding all the root vectors $`\rho ^a`$ of $`U(sn)`$ (or linear combinations of these) that satisfy
$$\rho ^aV_{g_1}=0\mathrm{mod}\text{/}\text{\_Z},\mathrm{\Pi }_{g_2}(\rho ^a)=\rho ^a.$$
(C.15)
If we choose twist vectors $`v=\frac{1}{N}(1,1,0)`$ resp. $`w=\frac{1}{M}(0,1,1)`$ for the action of $`g_1`$ resp. $`g_2`$ on $`\mathrm{C}\text{ }^3`$, then the matter representations are obtained from the root vectors $`\rho ^a`$ of $`U(sn)`$ (or linear combinations of these) that satisfy
$$\begin{array}{ccc}\rho ^aV_{g_1}=\frac{1}{N}\mathrm{mod}\text{/}\text{\_Z},\hfill & \mathrm{\Pi }_{g_2}(\rho ^a)=\rho ^a\hfill & \text{(}1^{\mathrm{st}}\text{ complex plane)\hspace{1em}or}\hfill \\ \rho ^aV_{g_1}=\frac{1}{N}\mathrm{mod}\text{/}\text{\_Z},\hfill & \mathrm{\Pi }_{g_2}(\rho ^a)=e^{2\pi i/M}\rho ^a\hfill & \text{(}2^{\mathrm{nd}}\text{ complex plane)\hspace{1em}or}\hfill \\ \rho ^aV_{g_1}=0\mathrm{mod}\text{/}\text{\_Z},\hfill & \mathrm{\Pi }_{g_2}(\rho ^a)=e^{2\pi i/M}\rho ^a\hfill & \text{(}3^{\mathrm{rd}}\text{ complex plane)}.\hfill \end{array}$$
(C.16)
Consider, for example, $`\mathrm{\Gamma }=\text{/}\text{\_Z}_N\times \text{/}\text{\_Z}_N`$ with discrete torsion $`ϵ=\omega _N`$, i.e. $`s=N`$ (these models were discussed in ). In this case, one has
$`V_{g_1}={\displaystyle \frac{1}{N}}(0^{n_0},1^{n_0},\mathrm{},(N1)^{n_0}),`$ (C.17)
$`\mathrm{\Pi }_{g_2}:(\stackrel{~}{\rho }_1^{(n_0)},\mathrm{},\stackrel{~}{\rho }_N^{(n_0)})(\stackrel{~}{\rho }_N^{(n_0)},\stackrel{~}{\rho }_1^{(n_0)},\mathrm{},\stackrel{~}{\rho }_{N1}^{(n_0)}).`$
The roots of $`U(Nn_0)`$ are of the form $`\rho ^a=(\underset{¯}{+,,0^{Nn_02}})`$, where underlining means that all permutations have to be considered. It is easy to see that there are $`n_0(n_01)`$ linear combinations of roots $`\rho _a`$ that satisfy (C.15):
$$(\underset{¯}{+,,0^{n_02}},0^{(N1)n_0})+(0^{n_0},\underset{¯}{+,,0^{n_02}},0^{(N2)n_0})+\mathrm{}+(0^{(N1)n_0},\underset{¯}{+,,0^{n_02}}).$$
Together with the $`n_0`$ Cartan generators they form the gauge group $`G=U(n_0)`$. Similarly, one finds $`3n_0^2n_0`$ linear combinations of roots $`\rho _a`$ (and $`n_0`$ Cartan generators) that satisfy (C.16):
$`1^{\mathrm{st}}\mathrm{plane}:`$ $`(\underset{¯}{,0^{n_01}},\underset{¯}{+,0^{n_01}},0^{(N2)n_0})+(0^{n_0},\underset{¯}{,0^{n_01}},\underset{¯}{+,0^{n_01}},0^{(N3)n_0})`$
$`+\mathrm{}+(\underset{¯}{+,0^{n_01}},0^{(N2)n_0},\underset{¯}{,0^{n_01}}),`$
$`2^{\mathrm{nd}}\mathrm{plane}:`$ $`(\underset{¯}{+,0^{n_01}},\underset{¯}{,0^{n_01}},0^{(N2)n_0})+\alpha ^1(0^{n_0},\underset{¯}{+,0^{n_01}},\underset{¯}{,0^{n_01}},0^{(N3)n_0})`$
$`+\mathrm{}+\alpha ^{(N1)}(\underset{¯}{,0^{n_01}},0^{(N2)n_0},\underset{¯}{+,0^{n_01}}),`$
$`3^{\mathrm{rd}}\mathrm{plane}:`$ $`(\underset{¯}{+,,0^{n_02}},0^{(N1)n_0})+\alpha (0^{n_0},\underset{¯}{+,,0^{n_02}},0^{(N2)n_0})`$
$`+\mathrm{}+\alpha ^{N1}(0^{(N1)n_0},\underset{¯}{+,,0^{n_02}}),`$
where $`\alpha =e^{2\pi i/N}`$. These are three matter fields in the adjoint representation of $`U(n_0)`$.
As a second example, consider $`\mathrm{\Gamma }=\text{/}\text{\_Z}_N\times \text{/}\text{\_Z}_2`$ with discrete torsion $`ϵ=1`$, i.e. $`s=2`$. In this case, one has
$`V_{g_1}={\displaystyle \frac{1}{N}}(0^{n_0},\left({\displaystyle \frac{N}{2}}\right)^{n_0},1^{n_1},\left({\displaystyle \frac{N}{2}}+1\right)^{n_1},\mathrm{},\left({\displaystyle \frac{N}{2}}1\right)^{n_{N/21}},(N1)^{n_{N/21}}),`$
$`\mathrm{\Pi }_{g_2}:(\stackrel{~}{\rho }_1^{(n_0)},\stackrel{~}{\rho }_2^{(n_0)},\stackrel{~}{\rho }_1^{(n_1)},\stackrel{~}{\rho }_2^{(n_1)},\mathrm{},\stackrel{~}{\rho }_1^{(n_{N/21})},\stackrel{~}{\rho }_2^{(n_{N/21})})`$ (C.18)
$`(\stackrel{~}{\rho }_2^{(n_0)},\stackrel{~}{\rho }_1^{(n_0)},\stackrel{~}{\rho }_2^{(n_1)},\stackrel{~}{\rho }_1^{(n_1)},\mathrm{},\stackrel{~}{\rho }_2^{(n_{N/21})},\stackrel{~}{\rho }_1^{(n_{N/21})}).`$
There are $`n_0(n_01)+n_1(n_11)+\mathrm{}+n_{N/21}(n_{N/21}1)`$ linear combinations of roots of $`U(2n)`$ that satisfy (C.15). Together with the $`n=_kn_k`$ Cartan generators, they form the gauge group $`G=U(n_0)\times U(n_1)\times \mathrm{}\times U(n_{N/21})`$. Similarly, one finds the following matter representations from the 3 complex planes:
$`1^{\mathrm{st}}\mathrm{plane}:`$ $`(\overline{\text{ }\text{ }\text{ }\text{ }\text{ }},\text{ }\text{ }\text{ }\text{ }\text{ },1\mathrm{I},\mathrm{},1\mathrm{I})+(1\mathrm{I},\overline{\text{ }\text{ }\text{ }\text{ }\text{ }},\text{ }\text{ }\text{ }\text{ }\text{ },1\mathrm{I},\mathrm{},1\mathrm{I})+\mathrm{}+(1\mathrm{I},\mathrm{},1\mathrm{I},\overline{\text{ }\text{ }\text{ }\text{ }\text{ }},\text{ }\text{ }\text{ }\text{ }\text{ })+(\text{ }\text{ }\text{ }\text{ }\text{ },1\mathrm{I},\mathrm{},1\mathrm{I},\overline{\text{ }\text{ }\text{ }\text{ }\text{ }})`$
$`2^{\mathrm{nd}}\mathrm{plane}:`$ $`(\text{ }\text{ }\text{ }\text{ }\text{ },\overline{\text{ }\text{ }\text{ }\text{ }\text{ }},1\mathrm{I},\mathrm{},1\mathrm{I})+(1\mathrm{I},\text{ }\text{ }\text{ }\text{ }\text{ },\overline{\text{ }\text{ }\text{ }\text{ }\text{ }},1\mathrm{I},\mathrm{},1\mathrm{I})+\mathrm{}+(1\mathrm{I},\mathrm{},1\mathrm{I},\text{ }\text{ }\text{ }\text{ }\text{ },\overline{\text{ }\text{ }\text{ }\text{ }\text{ }})+(\overline{\text{ }\text{ }\text{ }\text{ }\text{ }},1\mathrm{I},\mathrm{},1\mathrm{I},\text{ }\text{ }\text{ }\text{ }\text{ })`$
$`3^{\mathrm{rd}}\mathrm{plane}:`$ $`(adj,1\mathrm{I},\mathrm{},1\mathrm{I})+(1\mathrm{I},adj,1\mathrm{I},\mathrm{},1\mathrm{I})+\mathrm{}+(1\mathrm{I},\mathrm{},1\mathrm{I},adj)`$
It is instructive to see in a simple special case that the permutation $`\mathrm{\Pi }_{g_2}`$ in the shift formalism indeed corresponds to the action of $`\gamma _{g_2}`$ in the formalism using matrices. Let us take $`\mathrm{\Gamma }=\text{/}\text{\_Z}_4\times \text{/}\text{\_Z}_2`$. According to (C.6) the matrices $`\gamma _{g_1}`$, $`\gamma _{g_2}`$ are given by
$$\gamma _{g_1}=\left(\begin{array}{cccc}1\mathrm{I}_{n_0}& 0& 0& 0\\ 0& 1\mathrm{I}_{n_0}& 0& 0\\ 0& 0& i1\mathrm{I}_{n_1}& 0\\ 0& 0& 0& i1\mathrm{I}_{n_1}\end{array}\right),\gamma _{g_2}=\left(\begin{array}{cccc}0& 1\mathrm{I}_{n_0}& 0& 0\\ 1\mathrm{I}_{n_0}& 0& 0& 0\\ 0& 0& 0& 1\mathrm{I}_{n_1}\\ 0& 0& 1\mathrm{I}_{n_1}& 0\end{array}\right)$$
(C.19)
As $`\gamma _{g_1}`$ is diagonal, according to (C.2), it corresponds to a shift $`V_{g_1}=\frac{1}{4}(0^{n_0},2^{n_0},1^{n_1},3^{n_1})`$. The action of $`\gamma _{g_2}E_a\gamma _{g_2}^1`$ on the matrix $`E_a`$ is such that it permutes the associated roots $`\rho ^a`$ as $`(\stackrel{~}{\rho }_1^{(n_0)},\stackrel{~}{\rho }_2^{(n_0)},\stackrel{~}{\rho }_1^{(n_1)},\stackrel{~}{\rho }_2^{(n_1)})(\stackrel{~}{\rho }_2^{(n_0)},\stackrel{~}{\rho }_1^{(n_0)},\stackrel{~}{\rho }_2^{(n_1)},\stackrel{~}{\rho }_1^{(n_1)})`$. This coincides with the prescription given above. More precisely, $`\gamma _{g_2}E_a\gamma _{g_2}^1=\alpha E_a^{}`$, where $`E_a^{}`$ is the matrix corresponding to the permuted root vector. The phase $`\alpha =e^{m\pi i/M}`$, $`m=0,1,2,3`$, is related to an ambiguity in choosing a specific basis for the matrices $`E_a`$. It does not affect the orbifold spectrum and can therefore be ignored. However, in the orientifold case, it will be important to take this phase into account.
For a general $`\text{/}\text{\_Z}_N\times \text{/}\text{\_Z}_M`$ orbifold with discrete torsion $`ϵ`$, such that $`s<M`$, one has to take $`M/s`$ copies of the shift vector (C.13), with entries labelled by
$$(n_{0,0},n_{1,0},\mathrm{},n_{N/s1,0}|n_{0,1},n_{1,1},\mathrm{},n_{N/s1,1}|\mathrm{}|n_{0,M/s1},n_{1,M/s1},\mathrm{},n_{N/s1,M/s1}).$$
The permutation $`\mathrm{\Pi }_{g_2}`$ acts on each copy identically, as in (C). To represent the action of the phases $`\omega _M^l`$ that appear in (C.6) for $`s<M`$, one has to associate a shift vector also to the element $`g_2`$:
$$V_{g_2}=\frac{1}{M}(0^{\stackrel{~}{n}_0},1^{\stackrel{~}{n}_1},\mathrm{},\left(\frac{M}{s}1\right)^{\stackrel{~}{n}_{M/s1}}),\mathrm{with}\stackrel{~}{n}_l=\underset{k}{}n_{k,l}.$$
To determine the gauge group, one has to find all roots $`\rho ^a`$ of $`U(s_{k,l}n_{k,l})`$ that satisfy
$$\rho ^aV_{g_1}=0\mathrm{mod}\text{/}\text{\_Z},\rho ^aV_{g_2}=0\mathrm{mod}\text{/}\text{\_Z},\mathrm{\Pi }_{g_2}(\rho ^a)=\rho ^a.$$
(C.20)
The roots corresponding to the matter representations have to satisfy, for each complex plane, one of the following $`s`$ conditions:
$`\rho ^aV_{g_1}=v_i\mathrm{mod}\text{/}\text{\_Z},\rho ^aV_{g_2}=w_i{\displaystyle \frac{r}{s}}\mathrm{mod}\text{/}\text{\_Z},\mathrm{\Pi }_{g_2}(\rho ^a)=e^{2\pi ir/s}\rho ^a,`$
$`i=1,2,3,r=0,\mathrm{},s1,`$ (C.21)
$`v_1={\displaystyle \frac{1}{N}},v_2={\displaystyle \frac{1}{N}},v_3=0,w_1=0,w_2={\displaystyle \frac{1}{M}},w_3={\displaystyle \frac{1}{M}}.`$
As an example, let us compute the spectrum of the $`\text{/}\text{\_Z}_6\times \text{/}\text{\_Z}_6`$ orbifold with discrete torsion $`ϵ=e^{2\pi i/3}`$, i.e. $`s=3`$. The shift vectors and permutation are given by
$`V_{g_1}`$ $`=`$ $`{\displaystyle \frac{1}{6}}(0^{n_{0,0}},2^{n_{0,0}},4^{n_{0,0}},1^{n_{1,0}},3^{n_{1,0}},5^{n_{1,0}}|0^{n_{0,1}},2^{n_{0,1}},4^{n_{0,1}},1^{n_{1,1}},3^{n_{1,1}},5^{n_{1,1}}),`$
$`V_{g_2}`$ $`=`$ $`{\displaystyle \frac{1}{6}}(0^{\stackrel{~}{n}_0},1^{\stackrel{~}{n}_1}),`$ (C.22)
$`\mathrm{\Pi }_{g_2}:`$ $`(\stackrel{~}{\rho }_1^{n_{0,0}},\stackrel{~}{\rho }_2^{n_{0,0}},\stackrel{~}{\rho }_3^{n_{0,0}},\mathrm{},\stackrel{~}{\rho }_1^{n_{1,1}},\stackrel{~}{\rho }_2^{n_{1,1}},\stackrel{~}{\rho }_3^{n_{1,1}})(\stackrel{~}{\rho }_3^{n_{0,0}},\stackrel{~}{\rho }_1^{n_{0,0}},\stackrel{~}{\rho }_2^{n_{0,0}},\mathrm{},\stackrel{~}{\rho }_3^{n_{1,1}},\stackrel{~}{\rho }_1^{n_{1,1}},\stackrel{~}{\rho }_2^{n_{1,1}}).`$
There are $`n_{0,0}(n_{0,0}1)`$ linear combinations of roots,
$$(\underset{¯}{+,,0^{n_{0,0}2}},0,\mathrm{},0)+(0^{n_{0,0}},\underset{¯}{+,,0^{n_{0,0}2}},0,\mathrm{},0)+(0^{2n_{0,0}},\underset{¯}{+,,0^{n_{0,0}2}},0,\mathrm{},0),$$
that satisfy (C.20), and similarly for $`n_{1,0}`$, $`n_{0,1}`$, $`n_{1,1}`$. Together with the Cartan generators, they form the gauge group $`U(n_{0,0})\times U(n_{1,0})\times U(n_{0,1})\times U(n_{1,1})`$. Concerning the matter fields from the first plane, only $`r=0`$ is possible in (C). This gives the following representations:
$$(\text{ }\text{ }\text{ }\text{ }\text{ },\overline{\text{ }\text{ }\text{ }\text{ }\text{ }},1\mathrm{I},1\mathrm{I})+(\overline{\text{ }\text{ }\text{ }\text{ }\text{ }},\text{ }\text{ }\text{ }\text{ }\text{ },1\mathrm{I},1\mathrm{I})+(1\mathrm{I},1\mathrm{I},\text{ }\text{ }\text{ }\text{ }\text{ },\overline{\text{ }\text{ }\text{ }\text{ }\text{ }})+(1\mathrm{I},1\mathrm{I},\overline{\text{ }\text{ }\text{ }\text{ }\text{ }},\text{ }\text{ }\text{ }\text{ }\text{ }).$$
In the second plane, $`r=0`$ and $`r=1`$ are possible, yielding
$$r=0:(1\mathrm{I},\overline{\text{ }\text{ }\text{ }\text{ }\text{ }},\text{ }\text{ }\text{ }\text{ }\text{ },1\mathrm{I})+(\overline{\text{ }\text{ }\text{ }\text{ }\text{ }},1\mathrm{I},1\mathrm{I},\text{ }\text{ }\text{ }\text{ }\text{ }),r=1:(1\mathrm{I},\text{ }\text{ }\text{ }\text{ }\text{ },\overline{\text{ }\text{ }\text{ }\text{ }\text{ }},1\mathrm{I})+(\text{ }\text{ }\text{ }\text{ }\text{ },1\mathrm{I},1\mathrm{I},\overline{\text{ }\text{ }\text{ }\text{ }\text{ }}).$$
In the third plane, $`r=0`$ and $`r=2`$ are possible, yielding
$$r=0:(\text{ }\text{ }\text{ }\text{ }\text{ },1\mathrm{I},\overline{\text{ }\text{ }\text{ }\text{ }\text{ }},1\mathrm{I})+(1\mathrm{I},\text{ }\text{ }\text{ }\text{ }\text{ },1\mathrm{I},\overline{\text{ }\text{ }\text{ }\text{ }\text{ }}),r=2:(\overline{\text{ }\text{ }\text{ }\text{ }\text{ }},1\mathrm{I},\text{ }\text{ }\text{ }\text{ }\text{ },1\mathrm{I})+(1\mathrm{I},\overline{\text{ }\text{ }\text{ }\text{ }\text{ }},1\mathrm{I},\text{ }\text{ }\text{ }\text{ }\text{ }).$$
The open string spectrum of a general type IIB orientifold with or without discrete torsion is determined by the matrices $`\lambda `$ that satisfy (C.1) and, in addition,
$$\gamma _\mathrm{\Omega }\lambda ^{}\gamma _\mathrm{\Omega }^1=r(\mathrm{\Omega })\lambda ,$$
(C.23)
where $`\gamma _\mathrm{\Omega }`$ represents the action of the world-sheet parity $`\mathrm{\Omega }`$ on the Chan-Paton indices of the open string and $`r(\mathrm{\Omega })`$ the action of $`\mathrm{\Omega }`$ on the oscillator state of the open string. In complete analogy to the orbifold case, the spectrum can equivalently be obtained using the shift formalism. The additional condition (C.23) implies that the gauge group $`\stackrel{~}{G}`$ before the orbifold projection is $`SO(4_{k,l}n_{k,l})`$ resp. $`USp(4_{k,l}n_{k,l})`$ if $`\gamma _\mathrm{\Omega }`$ is symmetric resp. antisymmetric. In the case of orientifolds with discrete torsion, one has to find all root vectors $`\rho ^a`$ of $`\stackrel{~}{G}`$ that satisfy (C.20), (C). As a consequence of $`\mathrm{\Omega }g\mathrm{\Omega }=g`$, the matrices $`\gamma _g`$ now form a real projective representation and only discrete torsion $`ϵ=\pm 1`$ is possible.
In this case one must also distinguish between orientifolds with four different boundary conditions on the gauge bundle: $`\gamma _{g_1}^N=\pm 1\mathrm{I}`$, $`\gamma _{g_2}^M=\pm 1\mathrm{I}`$. In the shift formalism this translates to
$$NV_{g_1}=\{\begin{array}{cc}(1,\mathrm{},1)\mathrm{mod}\text{/}\text{\_Z}\hfill & \mathrm{if}\gamma _{g_1}^N=1\mathrm{I}\hfill \\ \frac{1}{2}(1,\mathrm{},1)\mathrm{mod}\text{/}\text{\_Z}\hfill & \mathrm{if}\gamma _{g_1}^N=1\mathrm{I}\hfill \end{array},$$
(C.24)
and similarly for $`MV_{g_2}`$. As mentioned above (below eq. (C.19)), there is an ambiguity in the choice of a basis for the Chan-Paton matrices $`\lambda `$ which leads to an additional phase in the action of $`g_2`$ on $`\lambda `$. It turns out that the above prescription, using a shift $`V_{g_2}`$ of the form (C) and requiring the conditions (C.20), (C), corresponds in the orientifold case to $`\gamma _{g_2}^M=1\mathrm{I}`$. However, using the shift $`V_{g_1}`$ of the form (C.13) for the group element $`g_1`$, corresponds to $`\gamma _{g_1}^N=1\mathrm{I}`$, as expected. If $`\gamma _{g_1}^N=1\mathrm{I}`$, $`\gamma _{g_2}^M=1\mathrm{I}`$, then the general form of the shifts given above, eqs. (C.13), (C), is modified to (using $`s=2`$ and taking $`M/2`$ copies of (C.13))
$`V_{g_1}`$ $`=`$ $`{\displaystyle \frac{1}{2N}}(1^{n_{0,0}},(N+1)^{n_{0,0}},3^{n_{1,0}},(N+3)^{n_{1,0}},\mathrm{},(N1)^{n_{N/21,0}},(2N1)^{n_{N/21,0}}|`$ (C.25)
$`1^{n_{0,1}},(N+1)^{n_{0,1}},\mathrm{}|\mathrm{}|\mathrm{}(2N1)^{n_{N/21,M/21}}),`$
$`V_{g_2}`$ $`=`$ $`{\displaystyle \frac{1}{2M}}(1^{\stackrel{~}{n}_0},3^{\stackrel{~}{n}_1},\mathrm{},(M1)^{\stackrel{~}{n}_{M/21}}),\mathrm{with}\stackrel{~}{n}_l={\displaystyle \underset{k}{}}n_{k,l}.`$
The permutation $`\mathrm{\Pi }_{g_2}`$ is not modified. The gauge group is determined by the roots $`\rho ^a`$ that satisfy one of the two conditions
$$\rho ^aV_{g_1}=0\mathrm{mod}\text{/}\text{\_Z},\rho ^aV_{g_2}=\frac{r}{2}\mathrm{mod}\text{/}\text{\_Z},\mathrm{\Pi }_{g_2}(\rho ^a)=(1)^r\rho ^a,r=0,1.$$
(C.26)
The matter fields are obtained from (C) by setting $`s=2`$. |
warning/0002/math0002153.html | ar5iv | text | # Dual group actions on C*–algebras and their description by Hilbert extensions
## 1 Introduction
In the Doplicher/Roberts theory (e.g. ) it is a central assumption that the center of the C\*–algebra $`𝒜`$ with which one starts the analysis is trivial, i.e. $`𝒵=𝒵(𝒜)=\mathrm{}\text{}`$ (in a more categorial notation the assumption reads $`(\iota ,\iota )=\mathrm{}\text{}`$, where $`\iota `$ denotes the unit object of the strict monoidal C\*–category, cf. ). From a systematical point of view it is interesting to study the properties and structural modifications of this theory if one assumes the presence of a nontrivial center $`𝒵\mathrm{}\text{}`$. For example, if $`(,\alpha _𝒢)`$ is a Hilbert C\*–system for a compact group $`𝒢`$ and if the corresponding fixed point algebra $`𝒜`$ has a nontrivial center that satisfies the relation $`𝒜^{}=𝒵`$, then the Galois correspondence does not hold anymore, i.e. we have the proper inclusion $`\alpha _𝒢\mathrm{stab}𝒜`$ in $`\mathrm{aut}`$ (cf. \[6, Section 7\]). Recall, that in the trivial center situation it is a fundamental result of the theory that $`\alpha _𝒢=\mathrm{stab}𝒜`$. As a further justification we can also mention that in other generalizations of the Doplicher/Roberts theory as well as in some applications in mathematical physics a nontrivial center plays, to a certain extent, a distinguished role .
In the present paper we continue the analysis of the presence of a nontrivial center in the construction of an extension algebra $``$ (cf. ). In particular, we study what we call dual group actions in the simple case where the group $`𝒳`$ is discrete and abelian (cf. with in the special case where $`𝒵=\mathrm{}\text{}`$). This investigations will be done at the abstract C\*–level which is the context of the Doplicher/Roberts theory mentioned above (cf. also ). On the other hand the results can be considered as a refinement of the study of twisted group algebras (twisted crossed products) on the concrete von Neumann algebra level (see e.g. ). For example, the decisive C\*–norm for the extension is defined intrisically and the natural representation (discussed e.g. by Sutherland) is related to the so–called regular representation that appears in the superselection theory . We hope that the present analysis will be useful to obtain a more general ‘inversion’ theorem, where endomorphisms of $`𝒜`$ are involved. Indeed, the main theorems in Section 3 suggest that for a more general inversion theory in the nontrivial center situation the cohomological aspects may be essential.
The paper is structured in 5 sections: in the following section we will introduce the notion of a Hilbert C\*–system and study some properties of the group homomorphism $`\mathrm{\Theta }:𝒳\mathrm{Out}𝒜`$. Hilbert C\*–systems are the result of the extension procedure mentioned above. In Section 3 we begin the study of the inverse (extension) problem: in particular it contains the result that a Hilbert extension exists iff there is a generalized 2–cocycle (to be defined there), and that in this case the set of all Hilbert extensions can be described in terms of the set of center–valued 2–cocycles of $`H^2(𝒳,𝒰(𝒵),\alpha _𝒳)`$ (cf. Theorems 3.4 and 3.8). In the next section we relate the previously obtained results to the special case of the Doplicher/Roberts frame, where $`𝒵=\mathrm{}\text{}`$. Finally, in Section 5 we give a representation of the Hilbert extension, which was already studied by Sutherland in the von Neumann case. In particular, we show that if there is a faithful state of $`𝒜`$, this representation coincides with the so–called regular representation that appears in superselection theory (cf. e.g. ) and the intrinsic C\*–norm turns out to be the operator norm of this representation.
## 2 Hilbert C\*–systems
A C\*-algebra $``$ together with a pointwise norm-continuous group homomorphism $`𝒢g\alpha _g\text{aut}`$ of a locally compact group $`𝒢`$ is called a C\*-system $`\{,\alpha _𝒢\}.`$ Let $`𝒜`$ be its fixed point algebra, i.e. $`𝒜:=\{A\alpha _gA=A,g𝒢\}`$. We denote by $`𝒜^c:=𝒜^{}`$ the relative commutant of $`𝒜`$ w.r.t. $``$. As is well-known, $`\alpha _g\text{}𝒜^c`$ is an automorphism of $`𝒜^c`$, so $`\{𝒜^c,\alpha _𝒢\}`$ is also a C\*-system. We call it the assigned C\*-system. The center $`𝒵(𝒜)`$ is denoted by $`𝒵`$.
In the following let $`𝒢`$ be compact and abelian so that $`\widehat{𝒢}=:𝒳`$ is abelian and discrete. The corresponding spectral projections w.r.t. $`\{,\alpha _𝒢\}`$ are denoted by $`\mathrm{\Pi }_\chi ,\chi 𝒳`$. Note that $`\mathrm{\Pi }_\iota =𝒜`$, where $`\iota `$ is the unit element of $`𝒳`$.
###### Definition 2.1
A C\*-system $`\{,\alpha _𝒢\},𝒢`$ compact abelian, is called a Hilbert C\*-system if spec$`\alpha _𝒢=𝒳`$ and if each spectral subspace $`\mathrm{\Pi }_\chi `$ contains a unitary $`U_\chi `$, i.e. $`𝒰(\mathrm{\Pi }_\chi )\mathrm{}`$.
If $`\{,\alpha _𝒢\}`$ is Hilbert, then $`\beta _\chi :=\text{ad}U_\chi \text{}𝒜`$ is an automorphism of $`𝒜`$, i.e. $`\beta _\chi \text{aut}𝒜.`$ We denote by $`\pi `$ the canonical homomorphism of aut$`𝒜`$ onto Out$`𝒜:=\text{aut}𝒜/\text{int}𝒜`$, where int$`𝒜`$ denotes the normal subgroup of all inner automorphisms of $`𝒜`$. Then
$$𝒳\chi \mathrm{\Theta }(\chi ):=\pi (\beta _\chi )\text{ Out}𝒜$$
(1)
is a group homomorphism of $`𝒳`$ into Out$`𝒜`$, i.e. we have
###### Lemma 2.2
To each Hilbert C\*–system $`\{,\alpha _𝒢\}`$, where $`𝒢`$ is compact abelian, there is canonically assigned a group homomorphism $`\mathrm{\Theta }:𝒳\text{Out}𝒜`$ given by (1).
Proof: Note that for $`\chi _1,\chi _2𝒳`$ we have that $`U_{\chi _1\chi _2}U_{\chi _2}^{}U_{\chi _1}^{}𝒜`$ and this implies that $`\beta _{\chi _1\chi _2}\beta _{\chi _2}^1\beta _{\chi _1}^1\mathrm{int}𝒜`$.
We mention next the characterization of those Hilbert C\*-systems where $`\mathrm{\Theta }`$ is an isomorphism and of those where the classes $`\mathrm{\Theta }(\chi )`$ are pairwise disjoint. Recall that $`\alpha ,\beta \text{aut}𝒜`$ are called disjoint if
$$(\alpha ,\beta ):=\{X𝒜X\alpha (A)=\beta (A)X\mathrm{for}\mathrm{all}A𝒜\}=0.$$
###### Proposition 2.3
* $`\mathrm{\Theta }`$ is a monomorphism iff no spectral subspace $`\mathrm{\Pi }_\chi 𝒜^c,\chi \iota `$, of the assigned C\*–system contains a unitary.
* The classes $`\mathrm{\Theta }(\chi )`$ are pairwise disjoint iff $`𝒜^c=𝒵`$, i.e. the relative commutant coincides with the center of $`𝒜`$.
Proof: For one of the directions of part (i) take a unitary $`U_\chi \mathrm{\Pi }_\chi (𝒜^c)`$ with $`\iota \chi 𝒳`$, so that the corresponding $`\beta _\chi =\mathrm{id}`$ and $`\pi (\beta _\chi )=\mathrm{int}𝒜`$. Thus $`\mathrm{\Theta }`$ is not injective. For the other implication take $`𝒳\chi _0\iota `$ with $`\chi _0\mathrm{ker}\mathrm{\Theta }`$, i.e. $`\mathrm{\Theta }(\chi _0)=\mathrm{int}𝒜`$. Thus there exists a unitary $`V𝒰(𝒜)`$ with $`\mathrm{ad}V=\mathrm{ad}U_{\chi _0}`$. From this we get $`V^{}U_{\chi _0}𝒰(𝒜^c)\mathrm{\Pi }_{\chi _0}()`$, i.e. $`\mathrm{\Pi }_{\chi _0}𝒜^c\mathrm{}`$.
Finally, part (ii) follows from \[7, Lemma 10.1.8\].
We mention several useful concepts for Hilbert C\*-systems $`\{,\alpha _𝒢\}`$ with a compact abelian group.
###### Definition 2.4
$`\beta \text{aut}𝒜`$ is called a canonical automorphism if $`\beta :=\text{ad}V\text{}𝒜,V_{\chi 𝒳}𝒰(\mathrm{\Pi }_\chi ).`$ The set of all canonical automorphisms is denoted by $`\mathrm{\Gamma }`$.
###### Remark 2.5
Note that for the set of canonical automorphisms we have $`\mathrm{int}𝒜\mathrm{\Gamma }\mathrm{aut}𝒜`$ and that for $`\alpha `$,$`\beta \mathrm{\Gamma }`$ the automorphisms $`\alpha \beta `$ and $`\beta \alpha `$ are unitarily equivalent. Furthermore, $`𝒳\mathrm{\Gamma }/\mathrm{int}𝒜`$ and the set $`\mathrm{\Gamma }`$ is sometimes called dual action on $`𝒜`$.
For any $`\gamma _1,\gamma _2\mathrm{\Gamma }`$ we write
$$\gamma _1\gamma _2\gamma _1^1\gamma _2^1=\text{ad}ϵ(\gamma _1,\gamma _2),$$
where $`ϵ(\gamma _1,\gamma _2)𝒰(𝒜)`$ and the class $`\widehat{ϵ}(\gamma _1,\gamma _2):=ϵ(\gamma _1,\gamma _2)\text{mod}𝒰(𝒵)`$ is uniquely defined.
###### Lemma 2.6
The permutators $`ϵ(,)`$ satisfy the following relations:
$`ϵ(\gamma _1,\gamma _2)ϵ(\gamma _2,\gamma _1)`$ $``$ $`\text{}\text{mod}𝒰(𝒵),\gamma _1,\gamma _2\mathrm{\Gamma },`$
$`ϵ(\iota ,\gamma )ϵ(\gamma ,\iota )`$ $``$ $`\text{}\text{mod}𝒰(𝒵),\gamma \mathrm{\Gamma },`$
$`\gamma _1(ϵ(\gamma _2,\gamma _3))ϵ(\gamma _1,\gamma _3)`$ $``$ $`ϵ(\gamma _1\gamma _2,\gamma _3)\text{mod}𝒰(𝒵),\gamma _1,\gamma _2,\gamma _3\mathrm{\Gamma },`$
$`A\gamma _1(B)ϵ(\gamma _1,\gamma _2)`$ $``$ $`ϵ(\gamma _1^{},\gamma _2^{})B\gamma _2(A)\text{mod}𝒰(𝒵),\gamma _1,\gamma _2,\gamma _1^{},\gamma _2^{}\mathrm{\Gamma }\mathrm{and}`$
$`A(\gamma _1,\gamma _1^{})𝒰(𝒜),B(\gamma _2,\gamma _2^{})𝒰(𝒜).`$
Proof: The first and second equations above are obvious. To prove the third one consider the the inner automorphism characterized by the l.h.s. of the equation:
$`\mathrm{ad}\left(\gamma _1(ϵ(\gamma _2,\gamma _3))ϵ(\gamma _1,\gamma _3)\right)`$ $`=`$ $`\mathrm{ad}\left(\gamma _1(ϵ(\gamma _2,\gamma _3))\right)\mathrm{ad}\left(ϵ(\gamma _1,\gamma _3)\right)`$
$`=`$ $`\gamma _1\mathrm{ad}(ϵ(\gamma _2,\gamma _3))\gamma _1^1\mathrm{ad}(ϵ(\gamma _1,\gamma _3))`$
$`=`$ $`\gamma _1(\gamma _2\gamma _3\gamma _2^1\gamma _3^1)\gamma _1^1(\gamma _1\gamma _3\gamma _1^1\gamma _3^1)=(\gamma _1\gamma _2)\gamma _3(\gamma _1\gamma _2)^1\gamma _3^1`$
$`=`$ $`\mathrm{ad}\left(ϵ(\gamma _1\gamma _2,\gamma _3)\right),`$
and this shows the desired relation. Finally, to prove the last equation recall that from the assumptions we have $`\gamma _1^{}=\mathrm{ad}(A)\gamma _1`$ and $`\gamma _2^{}=\mathrm{ad}(B)\gamma _2`$. From this we compute
$`\mathrm{ad}(ϵ(\gamma _1^{},\gamma _2^{}))`$ $`=`$ $`(\mathrm{ad}(A)\gamma _1)(\mathrm{ad}(B)\gamma _2)(\mathrm{ad}(A)\gamma _1)^1(\mathrm{ad}(B)\gamma _2)^1`$
$`=`$ $`\mathrm{ad}(A)\mathrm{ad}(\gamma _1(B))\underset{\mathrm{ad}(ϵ(\gamma _1,\gamma _2))}{\underset{}{\gamma _1\gamma _2\gamma _1^1\gamma _2^1}}\mathrm{ad}(\gamma _2(A))^1\mathrm{ad}(B)^1.`$
Therefore we get
$$\mathrm{ad}\left(ϵ(\gamma _1^{},\gamma _2^{})B\gamma _2(A)\right)=\mathrm{ad}\left(A\gamma _1(B)ϵ(\gamma _1,\gamma _2)\right)$$
which implies the last equation of the statement.
###### Definition 2.7
Let $`\beta _\chi \mathrm{\Theta }(\chi ),\chi 𝒳`$, with $`\beta _\iota =\text{id}_𝒜,`$ be a system of representatives, i.e. $`\pi (\beta _\chi )=\mathrm{\Theta }(\chi )`$. Then $`\beta _𝒳`$ is called a lifting of $`\mathrm{\Theta }`$ if $`𝒳\chi \beta _\chi \text{aut}𝒜`$ is a homomorphism.
###### Remark 2.8
For the notion of lifting see for example Jones . Sutherland says that $`\mathrm{\Theta }`$ splits if there is a lifting of $`\mathrm{\Theta }`$. If $`\mathrm{\Theta }`$ is an isomorphism then a lifting is also called monomorphic section (this latter name is used by Doplicher/Haag/Roberts ).
Results on the existence of liftings when $`𝒜`$ is a von Neumann algebra and in a more general context w.r.t. the group $`𝒳`$ (theory of Q-kernels) are due to Sutherland . Further, recall also the result of Doplicher/Haag/Roberts in the “automorphism case” of the superselection theory, where $`𝒵=\mathrm{}\text{}`$ and $`𝒜`$ is a so-called quasilocal algebra w.r.t. a net of local von Neumann algebras (see also ).
## 3 Hilbert extensions
The question concerning the description of $`\{,\alpha _𝒢\}`$ by $`𝒜`$ and ‘something else’ is called the reconstruction problem. It is posed, for example, by Takesaki \[23, p. 202\] and by Bratteli/Robinson \[8, p. 137\]. Also the superselection structures in algebraic quantum field theory are connected with the reconstruction problem (for the automorphism case see Doplicher/Haag/Roberts ).
From Lemma 2.2 it seems natural to consider the corresponding inverse problem, which is an extension problem. This is just the emphasis in the mentioned papers by Sutherland and Jones (see also Nakamura/Takeda ) as well as an essential aspect of the superselection theory (cf. ).
###### Definition 3.1
Let a system $`\{𝒜,\mathrm{\Theta }(𝒳)\}`$ be given where $`𝒳`$ is a discrete abelian group and where $`\mathrm{\Theta }:𝒳\text{Out}𝒜`$ is a homomorphism and put $`𝒢:=\widehat{𝒳}`$. A Hilbert C\*-system $`\{,\alpha _𝒢\}`$ is called a Hilbert extension of $`\{𝒜,\mathrm{\Theta }(𝒳)\}`$ if $`𝒜=\mathrm{\Pi }_\iota `$ and $`\mathrm{\Theta }(𝒳)`$ coincides with the homomorphism given by Lemma 2.2.
Now let $`\{𝒜,\mathrm{\Theta }(𝒳)\}`$ and $`𝒢`$ be given as in the previous definition. As it is pointed out, for example in , a crucial object for the extension problem is the so-called obstruction Ob$`\mathrm{\Theta }`$. We recall the relevant relations: Choose a system $`\beta _\chi \mathrm{\Theta }(\chi ),\chi 𝒳,\beta _\iota :=\text{id}_𝒜`$ of representatives. Then
$$\beta _{\chi _1}\beta _{\chi _2}=\text{ad}(\omega (\chi _1,\chi _2))\beta _{\chi _1\chi _2},$$
(2)
where
$$𝒳\times 𝒳(\chi _1,\chi _2)\omega (\chi _1,\chi _2)𝒰(𝒜)$$
(3)
and we have the intertwining property
$$\omega (\chi _1,\chi _2)(\beta _{\chi _1\chi _2},\beta _{\chi _1}\beta _{\chi _2}),$$
(4)
which is implied by (2). Moreover we have
$$\omega (\iota ,\chi )=\omega (\chi ,\iota )=\text{}.$$
(5)
Now associativity yields
$$\text{ad}(\omega (\chi _1,\chi _2)\omega (\chi _1\chi _2,\chi _3))=\text{ad}(\beta _{\chi _1}(\omega (\chi _2,\chi _3))\omega (\chi _1,\chi _2\chi _3))$$
so that there is $`\gamma (\chi _1,\chi _2,\chi _3)𝒰(𝒵)`$ with
$$\omega (\chi _1,\chi _2)\omega (\chi _1\chi _2,\chi _3)=\gamma (\chi _1,\chi _2,\chi _3)\beta _{\chi _1}(\omega (\chi _2,\chi _3))\omega (\chi _1,\chi _2\chi _3).$$
If $`\gamma (\chi _1,\chi _2,\chi _3)=\text{}`$ for all $`\chi _1,\chi _2,\chi _3𝒳`$ we obtain the equation
$$\omega (\chi _1,\chi _2)\omega (\chi _1\chi _2,\chi _3)=\beta _{\chi _1}(\omega (\chi _2,\chi _3))\omega (\chi _1,\chi _2\chi _3).$$
(6)
Obviously, the existence of a system of representatives $`\beta _𝒳`$ such that equation (6) has a solution $`\omega `$ equipped with the properties (3)–(5) is necessary for the existence of a Hilbert extension. Even more, the existence of such a solution is also sufficient for the existence of a Hilbert extension.
###### Definition 3.2
A function $`\omega `$, assigned to a given system $`\beta _𝒳`$ of representatives of $`\mathrm{\Theta }(𝒳)`$, equipped with the properties (3)–(6) is called a generalized 2-cocycle.
One calculates easily that the existence of a generalized 2-cocycle is independent of the choice of the system $`\beta _𝒳`$ of representatives. Further, a generalized cocycle $`\omega `$ for $`\beta _𝒳`$ satisfies the relation
$$\text{ad}(\omega (\chi _1,\chi _2)\omega (\chi _2,\chi _1)^1)=\beta _{\chi _1}\beta _{\chi _2}\beta _{\chi _1}^1\beta _{\chi _2}^1.$$
The existence of a lifting of $`\mathrm{\Theta }`$ can be expressed in terms of generalized 2-cocycles as follows.
###### Lemma 3.3
There exists a lifting $`\beta _𝒳`$ of $`\mathrm{\Theta }`$ iff to each system $`\gamma _𝒳`$ of representatives there corresponds a generalized 2–cocycle $`\omega `$ of the form
$$\omega (\chi _1,\chi _2)\gamma _{\chi _1}(V_{\chi _2}^1)V_{\chi _1}^1V_{\chi _1\chi _2}\mathrm{mod}𝒰(𝒵),$$
where $`V_\chi 𝒰(𝒜),V_\iota =\text{}.`$ In this case, i.e. if there is a lifting $`\beta _𝒳`$, then a corresponding generalized 2–cocycle $`\omega `$ is given by $`\omega (\chi _1,\chi _2)=\text{}`$ for all $`\chi _1,\chi _2𝒳.`$
Proof: Let $`\beta _\chi =\mathrm{ad}(V_\chi )\gamma _\chi `$, $`V_\chi 𝒰(𝒜)`$, $`\chi 𝒳`$. Now if $`\omega (\chi _1,\chi _2)=\gamma _{\chi _1}(V_{\chi _2}^1)V_{\chi _1}^1V_{\chi _1\chi _2}Z`$ for some $`Z𝒰(𝒵)`$, then we have on the one hand $`\beta _{\chi _1\chi _2}=\mathrm{ad}(V_{\chi _1\chi _2})\gamma _{\chi _1\chi _2}`$ and on the other
$$\beta _{\chi _1}\beta _{\chi _2}=(\mathrm{ad}(V_{\chi _1})\gamma _{\chi _1})(\mathrm{ad}(V_{\chi _2})\gamma _{\chi _2})=\mathrm{ad}\left(V_{\chi _1}\gamma _{\chi _1}(V_{\chi _2})\omega (\chi _1,\chi _2)\right)\gamma _{\chi _1\chi _2},$$
which using the assumption on $`\omega `$ and the fact that $`\mathrm{ad}(V_{\chi _1\chi _2}Z)=\mathrm{ad}(V_{\chi _1\chi _2})`$, implies that $`\beta _{\chi _1\chi _2}=\beta _{\chi _1}\beta _{\chi _2}`$, i.e. there is a lift of $`\mathrm{\Theta }`$. To prove the converse let $`\beta _{\chi _1\chi _2}=\beta _{\chi _1}\beta _{\chi _2}`$, so that from the above relations we have
$$\mathrm{ad}(V_{\chi _1\chi _2})=\mathrm{ad}\left(V_{\chi _1}\gamma _{\chi _1}(V_{\chi _2})\omega (\chi _1,\chi _2)\right),$$
which implies $`\omega (\chi _1,\chi _2)=\gamma _{\chi _1}(V_{\chi _2}^1)V_{\chi _1}^1V_{\chi _1\chi _2}\mathrm{mod}𝒰(𝒵)`$.
###### Theorem 3.4
Let $`\omega `$ be a generalized 2–cocycle for the system $`\beta _𝒳`$ of representatives. Then there is a Hilbert extension $`\{,\alpha _𝒢\}`$ of $`\{𝒜,\mathrm{\Theta }(𝒳)\}.`$
Proof: The proof consists of several steps that correspond to gradually imposing a richer structure on an initially considered $`𝒜`$–left module:
1. Indeed, choose first system of 1-dimensional linear spaces, generated by abstract elements $`U_\chi `$, $`\chi 𝒳`$, $`U_\iota :=\text{}𝒜`$. Form the $`𝒜`$–left modules $`𝒜\mathrm{}\mathrm{𝕌}_\chi `$ and $`_0:=_\chi (𝒜\mathrm{}\mathrm{𝕌}_\chi \mathrm{}\mathrm{}`$ By identification $`A\text{}A,\text{}U_\chi U_\chi `$ one has
$$_0=\{\underset{\chi \text{, finite sum}}{}A_\chi U_\chi A_\chi 𝒜\},$$
where $`\{U_\chi \chi 𝒳\}`$ forms an abstract $`𝒜`$–module basis.
2. Next we want to equip $`_0`$ with a multiplication structure. First $`_0`$ becomes an $`𝒜`$–bimodule extending linearly the following definition
$$U_\chi A:=\beta _\chi (A)U_\chi ,A𝒜,\chi 𝒳,$$
where $`\beta _𝒳`$ is the system of representatives to which we associate the generalized cocycle $`\omega `$. Now the product structure is finally specified by putting
$$U_{\chi _1}U_{\chi _2}:=\omega (\chi _1,\chi _2)U_{\chi _1\chi _2},\chi _1,\chi _2𝒳,$$
where the cocycle equation (6) guarantees that the product is associative and the boundary conditions (5) lead to $`U_\chi \text{}=\text{}U_\chi =U_\chi `$. Note that the preceding product structure already implies that the $`U_\chi `$ are invertible. Indeed, it can be checked easily that the inverse is given explicitly by
$$U_\chi ^1:=\beta _{\chi ^1}\left(\omega (\chi ,\chi ^1)^1\right)U_{\chi ^1}$$
(use for example the relation $`\beta _\chi (\omega (\chi ^1,\chi ))=\omega (\chi ,\chi ^1)`$, which follows from the cocycle equation (6) by putting $`\chi _1:=\chi `$, $`\chi _2:=\chi ^1`$ and $`\chi _3=\chi `$).
3. The following step consists in defining a \*–structure on $`_0`$. This is done by putting
$$U_\chi ^{}:=\omega (\chi ^1,\chi )^{}U_{\chi ^1}\mathrm{and}(AU_\chi )^{}:=U_\chi ^{}A^{}.$$
We still have to check that this definition is consistent, in particular with the product structure in $`_0`$, i.e. we have to verify:
$$(U_\chi ^{})^{}=U_\chi ,(U_\chi A)^{}=A^{}U_\chi ^{}\mathrm{and}(U_{\chi _1}U_{\chi _2})^{}=U_{\chi _2}^{}U_{\chi _1}^{}.$$
(7)
For the first equation we have
$`(U_\chi ^{})^{}`$ $`=`$ $`\left(\omega (\chi ^1,\chi )^{}U_{\chi ^1}\right)^{}=U_{\chi ^1}^{}\omega (\chi ^1,\chi )=\omega (\chi ,\chi ^1)^{}U_\chi \omega (\chi ^1,\chi )`$
$`=`$ $`\omega (\chi ,\chi ^1)^{}\beta _\chi \left(\omega (\chi ^1,\chi )^{}\right)U_\chi =\omega (\chi ,\chi ^1)^{}\omega (\chi ,\chi ^1)U_\chi `$
$`=`$ $`U_\chi `$
The second equation in (7) can also be checked immediately from the definitions considered above. For the last equation we will consider the two sides separately: for the r.h.s. we have
$`U_{\chi _2}^{}U_{\chi _1}^{}`$ $`=`$ $`\omega (\chi _2^1,\chi _2)^{}U_{\chi _2^1}\omega (\chi _1^1,\chi _1)^{}U_{\chi _1^1}`$
$`=`$ $`\omega (\chi _2^1,\chi _2)^{}\beta _{\chi _2^1}\left(\omega (\chi _1^1,\chi _1)^{}\right)U_{\chi _2^1}U_{\chi _1^1}`$
$`=`$ $`\omega (\chi _2^1,\chi _2)^{}\beta _{\chi _2^1}\left(\omega (\chi _1^1,\chi _1)^{}\right)\omega (\chi _2^1,\chi _1^1)U_{(\chi _1\chi _2)^1}`$
$`=`$ $`\omega (\chi _2^1,\chi _2)^{}\omega ((\chi _1\chi _2)^1,\chi _1)^{}\underset{\text{}}{\underset{}{\omega (\chi _2^1,\chi _1^1)^{}\omega (\chi _2^1,\chi _1^1)}}U_{(\chi _1\chi _2)^1},`$
where we have used the relation
$$\beta _{\chi _2^1}(\omega (\chi _1^1,\chi _1))=\omega (\chi _2^1,\chi _1^1)\omega (\chi _2^1\chi _1^1,\chi _1),$$
which again follows from the cocycle equation (6) taking now $`\chi _1:=\chi _2^1`$, $`\chi _2:=\chi _1^1`$ and $`\chi _3=\chi _1`$. Now the l.h.s. reads
$`(U_{\chi _1}U_{\chi _2})^{}`$ $`=`$ $`U_{\chi _1\chi _2}^{}\omega (\chi _1,\chi _2)^{}=\omega ((\chi _1\chi _2)^1,\chi _1\chi _2)^{}U_{(\chi _1\chi _2)^1}\omega (\chi _1,\chi _2)^{}`$
$`=`$ $`\omega ((\chi _1\chi _2)^1,\chi _1\chi _2)^{}\beta _{(\chi _1\chi _2)^1}\left(\omega (\chi _1,\chi _2)^{}\right)U_{(\chi _1\chi _2)^1}.`$
Thus to show the last equation in (7) we need to prove that
$$\omega ((\chi _1\chi _2)^1,\chi _1\chi _2)^{}\beta _{(\chi _1\chi _2)^1}\left(\omega (\chi _1,\chi _2)^{}\right)=\omega (\chi _2^1,\chi _2)^{}\omega ((\chi _1\chi _2)^1,\chi _1)^{}$$
or taking adjoints
$$\beta _{(\chi _1\chi _2)^1}\left(\omega (\chi _1,\chi _2)\right)\omega ((\chi _1\chi _2)^1,\chi _1\chi _2)=\omega ((\chi _1\chi _2)^1,\chi _1)\omega (\chi _2^1,\chi _2).$$
But the preceding equation is nothing else than the cocycle equation (7) with $`\chi _1:=(\chi _1\chi _2)^1`$, $`\chi _2:=\chi _1`$ and $`\chi _3:=\chi _2`$. Finally, note that since $`\beta _{\chi ^1}\left(\omega (\chi ,\chi ^1)^1\right)=\omega (\chi ^1,\chi )^{}`$ we also have that the $`U_\chi `$, are unitary, i.e. $`U_\chi ^{}=U_\chi ^1`$, $`\chi 𝒳`$.
4. Here we will define a representation of the compact abelian group $`𝒢=\widehat{𝒳}`$ in terms of automorphisms of the \*–algebra $`_0`$. The automorphisms are fixed by putting
$$\alpha _g(U_\chi ):=\chi (g)U_\chi \mathrm{and}\alpha _g(AU_\chi ):=A\alpha _g(U_\chi )=\chi (g)AU_\chi ,g𝒢,A𝒜,\chi 𝒳.$$
First we check that with the definition above the $`\alpha _g`$ is indeed an automorphism compatible with the structure in $`_0`$:
$`\alpha _g\left(U_{\chi _1}U_{\chi _2}\right)`$ $`=`$ $`\alpha _g\left(\omega (\chi _1,\chi _2)U_{\chi _1\chi _2}\right)=(\chi _1\chi _2)(g)\omega (\chi _1,\chi _2)U_{\chi _1\chi _2}`$
$`=`$ $`\chi _1(g)\chi _2(g)U_{\chi _1}U_{\chi _2}=\alpha _g\left(U_{\chi _1}\right)\alpha _g\left(U_{\chi _2}\right)`$
and
$`\alpha _g\left(U_\chi ^{}\right)`$ $`=`$ $`\alpha _g\left(\omega (\chi ^1,\chi )^{}U_{\chi ^1}\right)=(\chi ^1)(g)\omega (\chi ^1,\chi )^{}U_{\chi ^1}`$
$`=`$ $`\overline{\chi }(g)U_\chi ^{}=\alpha _g\left(U_\chi \right)^{}.`$
It can be also easily seen that the assignment $`𝒢g\alpha _g\mathrm{aut}_0`$ is an injective group homomorphism. Finally, note that the fixed point algebra of the previous action coincides with $`𝒜`$, i.e. for $`F_0`$, $`\alpha _g(F)=F`$ for all $`g𝒢`$ iff $`F𝒜`$. Indeed, for an arbitrary element $`_\chi A_\chi U_\chi _0`$ the equation $`_\chi \chi (g)A_\chi U_\chi =_\chi A_\chi U_\chi `$, $`g𝒢`$, implies by the base property of the $`U_\chi `$ that $`\chi (g)A_\chi =A_\chi `$ , $`g𝒢`$, $`\chi 𝒳`$. Therefore if $`\chi _0\iota `$, then there is a $`g_0𝒢`$ with $`\chi _0(g_0)1`$ and this shows that $`A_{\chi _0}=0`$. The converse implication is obvious.
5. Finally, to specify a C\*–norm on $`_0`$ we introduce the following $`𝒜`$–valued scalar product (note the variation w.r.t. the definition in \[2, p. 101\]):
$$F_1,F_2:=\underset{\chi }{}\beta _\chi ^1(A_\chi ^{}B_\chi ),\mathrm{where}F_1=\underset{\chi }{}A_\chi U_\chi ,F_2=\underset{\chi }{}B_\chi U_\chi _0.$$
This scalar product satisfies the properties
$$F_1,F_2^{}=F_2,F_1,F_1,F_10\mathrm{and}F_1,F_1=0\mathrm{iff}F_1=0.$$
Next we show that
$$F_1,F_2=\mathrm{\Pi }_\iota (F_1^{}F_2),$$
Indeed, using the definitions above we have
$$F_1^{}F_2=\underset{\chi _1,\chi _2}{}U_{\chi _1}^{}A_{\chi _1}^{}B_{\chi _2}U_{\chi _2}=\underset{\chi _1,\chi _2}{}\omega (\chi _1^1,\chi _1)^{}\beta _{\chi _1^1}\left(A_{\chi _1}^{}B_{\chi _2}\right)\omega (\chi _1^1,\chi _2)U_{\chi _1^1\chi _2}.$$
Putting, $`\chi _1=\chi _2=\chi `$ in the preceding expression we get
$`\mathrm{\Pi }_\iota (F_1^{}F_2)`$ $`=`$ $`{\displaystyle \underset{\chi }{}}\omega (\chi ^1,\chi )^{}\beta _{\chi ^1}\left(A_\chi ^{}B_\chi \right)\omega (\chi ^1,\chi )`$
$`=`$ $`{\displaystyle \underset{\chi }{}}\omega (\chi ^1,\chi )^{}\omega (\chi ^1,\chi )\beta _\chi ^1\left(A_\chi ^{}B_\chi \right)\omega (\chi ^1,\chi )^{}\omega (\chi ^1,\chi )`$
$`=`$ $`F_1,F_2,`$
where for the second equation before we have used eq. (2) in the form $`\beta _{\chi ^1}=\mathrm{ad}(\omega (\chi ^1,\chi ))\beta _\chi ^1`$. In particular the relation above implies the following invariance property: $`\alpha _g(F_1),\alpha _g(F_2)=F_1,F_2`$, $`g𝒢`$.
Define next the following norm on $`_0`$ by
$$|F|:=F,F^{\frac{1}{2}},F_0,$$
and the representation of $`_0`$ on $`(_0,||)`$ in terms of multiplication operators
$$\rho (F)X:=FX,F,X_0.$$
Note that by the definition of the $`𝒜`$–valued scalar product the property $`\rho (F^{})=\rho (F)^{}`$, $`F_0`$, holds. Now using the corresponding operator norm we introduce
$$F_{}:=|\rho (F)|_{op},F_0,$$
which by similar arguments as in \[2, p. 102-103\] satisfies the C\*–property $`F^{}F_{}=F_{}^2`$. Further, it satisfies also (cf. again the previous reference)
$$A_{}=A,A𝒜\mathrm{and}\alpha _g(F)_{}=F_{},g𝒢,F_0.$$
Therefore, we can finally extend $`\alpha _g`$ isometrically from $`_0`$ to
$$:=\mathrm{clo}_{_{}}(_0).$$
Further, $`\alpha _𝒢\mathrm{aut}`$ is norm continuous w.r.t. the pointwise norm convergence, because for any $`F_0=_\chi A_\chi U_\chi _0`$ we have
$$\alpha _{g_1}(F_0)\alpha _{g_2}(F_0)_{}=\underset{\chi }{}\left(\chi (g_1)\chi (g_2)\right)A_\chi U_\chi _{}\underset{\chi }{}|\chi (g_1)\chi (g_2)|A_\chi .$$
By construction we also have that $`U_\chi \mathrm{\Pi }_\chi ()`$, $`\chi 𝒳`$. Therefore from the definitions of Sections 2 and 3 we have constructed a Hilbert C\*–extension $`\{,\alpha _𝒢\}`$ of $`\{𝒜,\mathrm{\Gamma }\}`$ and the proof is concluded.
Using now Lemma 3.3 one has
###### Corollary 3.5
If there is a lifting of $`\mathrm{\Theta }`$, then there is a Hilbert extension of $`\{𝒜,\mathrm{\Theta }(𝒳)\}`$, corresponding to $`\omega =\text{}`$.
###### Remark 3.6
The construction in the proof of the previous theorem generalizes to the nontrivial center situation the procedure already presented (with small modifications) in \[2, Section 3.6\].
The second problem consists in the description of all Hilbert extensions. For this purpose let $`\mathrm{\Omega }(𝒳,𝒰(𝒵),\beta _𝒳)`$ be the set of all $`𝒰(𝒵)`$–valued 2–cocycles $`\lambda `$, i.e. $`\lambda `$ satisfies equation (6) and condition (5), but (3),(4) are replaced by $`\lambda (\chi _1,\chi _2)𝒰(𝒵)`$. For example, $`\lambda (\chi _1,\chi _2):=\text{}`$ for all $`\chi _1,\chi _2𝒳`$ is such a cocycle. Further let $`\mathrm{\Omega }_0(𝒳,𝒰(𝒵),\beta _𝒳)`$ be the set of all $`𝒰(𝒵)`$–valued coboundaries $`Z`$, i.e.
$$Z(\chi _1,\chi _2):=\frac{Z(\chi _1)\beta _{\chi _1}(Z(\chi _2))}{Z(\chi _1\chi _2)},$$
where $`Z()`$ is a $`𝒰(𝒵)`$-valued 1-cycle, $`Z(\iota )=\text{}`$. Then $`Z`$ is a $`𝒰(𝒵)`$-valued 2-cocycle, $`\mathrm{\Omega }\mathrm{\Omega }_0`$. As usual, $`\mathrm{\Omega }`$ and $`\mathrm{\Omega }_0`$ are abelian groups w.r.t. pointwise multiplication and the second cohomology is given by $`H^2(𝒳,𝒰(𝒵),\beta _𝒳):=\mathrm{\Omega }/\mathrm{\Omega }_0.`$
Next we need the concept of $`𝒜`$module isomorphism of Hilbert extensions.
###### Definition 3.7
Let $`\{^1,\alpha _𝒢^1\},\{^2,\alpha _𝒢^2\}`$ be Hilbert extensions of $`\{𝒜,\mathrm{\Theta }(𝒳)\}`$. They are called $`𝒜`$–module isomorphic if there is an algebraic isomorphism $`\mathrm{\Phi }:^1^2`$, with $`\mathrm{\Phi }(A)=A`$ for all $`A𝒜`$ and $`\mathrm{\Phi }\alpha _g^1=\alpha _g^2\mathrm{\Phi }`$ for all $`g𝒢`$.
###### Theorem 3.8
Let $`\omega _0`$ be a generalized 2–cocycle. Then:
* Each $`𝒰(𝒵)`$–valued 2–cocycle $`\lambda `$ yields a Hilbert extension generated by the generalized 2–cocycle $`\omega :=\lambda \omega _0`$ and each Hilbert extension is generated by some $`𝒰(𝒵)`$–valued 2–cocycle $`\lambda `$ via $`\omega :=\lambda \omega _0.`$
* Two Hilbert extensions are $`𝒜`$–module isomorphic iff the generating generalized 2–cocycles $`\omega _1,\omega _2`$ differ only by a $`𝒰(𝒵)`$–valued coboundary $`Z`$, i.e. $`\omega _1=Z\omega _2.`$
Proof: (i) If two generalized 2–cocycles $`\omega _1,\omega _2`$ are given, then note first that $`\lambda (\chi _1,\chi _2):=\omega _1(\chi _1,\chi _2)\omega _2(\chi _1,\chi _2)^1𝒰(𝒵)`$ for all $`\chi _1,\chi _2`$, because of condition (4). Further, eq. (5) follows from the corresponding properties of $`\omega _1`$ and $`\omega _2`$. Finally, the cocycle equation for $`\lambda (\chi _1,\chi _2)`$ is a consequence of the following computation:
$`\lambda (\chi _1,\chi _2)\lambda (\chi _1\chi _2,\chi _3)`$ $`=`$ $`\omega _1(\chi _1,\chi _2)\omega _2(\chi _1,\chi _2)^1\omega _1(\chi _1\chi _2,\chi _3)\omega _2(\chi _1\chi _2,\chi _3)^1`$
$`=`$ $`\omega _1(\chi _1,\chi _2)\omega _1(\chi _1\chi _2,\chi _3)\omega _2(\chi _1\chi _2,\chi _3)^1\omega _2(\chi _1,\chi _2)^1`$
$`=`$ $`(\omega _1(\chi _1,\chi _2)\omega _1(\chi _1\chi _2,\chi _3))(\omega _2(\chi _1,\chi _2)\omega _2(\chi _1\chi _2,\chi _3))^1`$
$`=`$ $`\beta _{\chi _1}(\omega _1(\chi _2,\chi _3))\omega _1(\chi _1,\chi _2\chi _3)(\beta _{\chi _1}(\omega _2(\chi _2,\chi _3))\omega _2(\chi _1,\chi _2\chi _3))^1`$
$`=`$ $`\beta _{\chi _1}(\omega _1(\chi _2,\chi _3))\omega _1(\chi _1,\chi _2\chi _3)\omega _2(\chi _1,\chi _2\chi _3)^1\beta _{\chi _1}(\omega _2(\chi _2,\chi _3))^1`$
$`=`$ $`\beta _{\chi _1}(\omega _1(\chi _2,\chi _3)\omega _2(\chi _2,\chi _3)^1)\omega _1(\chi _1,\chi _2\chi _3)\omega _2(\chi _1,\chi _2\chi _3)^1`$
$`=`$ $`\beta _{\chi _1}(\lambda (\chi _2,\chi _3))\lambda (\chi _1,\chi _2\chi _3),`$
i.e. if one fixes a generalized 2-cocycle $`\omega _0`$, then $`\omega :=\lambda \omega _0`$ runs through all generalized 2–cocycles $`\omega `$ if $`\lambda `$ runs through all $`𝒰(𝒵)`$–valued 2–cocycles in $`\mathrm{\Omega }(𝒳,𝒰(𝒵),\beta _𝒳)`$.
(ii) Let $`\{^1,\alpha _𝒢^1\}`$ and $`\{^2,\alpha _𝒢^2\}`$ be two Hilbert extensions of $`\{𝒜,\mathrm{\Theta }(𝒳)\}`$ and denote the corresponding set of abstract unitaries by $`\{U_\chi \chi 𝒳\}`$ resp. $`\{V_\chi \chi 𝒳\}`$.
Suppose first that there exists coboundary $`Z\mathrm{\Omega }_0(𝒳,𝒰(𝒵),\beta _𝒳)`$, where $`\beta _𝒳`$ is system of representatives in $`\mathrm{\Theta }`$, such that the corresponding generalized cocycles $`\omega _1`$ and $`\omega _2`$ satisfy $`\omega _1=Z\omega _2`$. In this case we will show that the extensions are isomorphic. Indeed, define the isomorphism by
$$\mathrm{\Phi }(AU_\chi ):=AZ(\chi )V_\chi ,A𝒜,\chi 𝒳,$$
and extend it by linearity to the corresponding left $`𝒜`$–module. Now $`\mathrm{\Phi }`$ is even a \*–homomorphism between the \*–algebras $`_0^1`$ and $`_0^2`$ that are defined in step 3 of the proof of Theorem 3.4. This follows from the following computations:
$`\mathrm{\Phi }(U_\chi A)`$ $`=`$ $`\mathrm{\Phi }\left(\beta _\chi (A)U_\chi \right)=Z(\chi )V_\chi A=\mathrm{\Phi }(U_\chi )\mathrm{\Phi }(A),`$
$`\mathrm{\Phi }(U_\chi U_\chi ^{})`$ $`=`$ $`\mathrm{\Phi }\left(\omega _1(\chi ,\chi ^{})U_{\chi \chi ^{}}\right)=Z(\chi ,\chi ^{})\omega _2(\chi ,\chi ^{})Z(\chi \chi ^{})V_{\chi \chi ^{}}`$
$`=`$ $`{\displaystyle \frac{Z(\chi )\beta _\chi (Z(\chi ^{}))}{Z(\chi \chi ^{})}}Z(\chi \chi ^{})V_\chi V_\chi ^{}=Z(\chi )V_\chi Z(\chi ^{})V_\chi ^{}=\mathrm{\Phi }(U_\chi )\mathrm{\Phi }(U_\chi ^{}),`$
$`\mathrm{\Phi }(U_\chi ^{})`$ $`=`$ $`\mathrm{\Phi }\left(\omega _1(\chi ^1,\chi )^{}U_{\chi ^1}\right)=Z(\chi ^1,\chi )^{}\omega _2(\chi ^1,\chi )^{}Z(\chi ^1)V_{\chi ^1}`$
$`=`$ $`Z(\chi ^1)^{}\beta _{\chi ^1}(Z(\chi ))^{}Z(\chi ^1)\omega _2(\chi ^1,\chi )^{}V_{\chi ^1}=(Z(\chi )V_\chi )^{}=\mathrm{\Phi }(U_\chi )^{},`$
where $`\chi ,\chi ^{}𝒳`$, $`A𝒜`$. Note further that on $`_0^1`$ we already have $`\mathrm{\Phi }\alpha _g^1=\alpha _g^2\mathrm{\Phi }`$, $`g𝒢`$, since for any $`\chi 𝒳`$ we have
$$\mathrm{\Phi }\alpha _g^1(AU_\chi )=\chi (g)AZ(\chi )V_\chi =\alpha _g^2(AZ(\chi )V_\chi )=\alpha _g^2\mathrm{\Phi }(AU_\chi ).$$
Recall that $`\mathrm{\Phi }`$ is a bijection between $`_0^1`$ and $`_0^2`$ and we will finish this part of the proof if we can also show that $`\mathrm{\Phi }`$ is even an isometry w.r.t the corresponding C\*–norms, because in this case we can isometrically extend $`\mathrm{\Phi }`$ to the desired Hilbert extension isomorphism $`\mathrm{\Phi }:^1^2`$. Now denote by $`,_k`$ the $`𝒜`$–valued scalar products on $`_0^k`$, $`k=1,2`$, given in step 5 of the proof of Theorem 3.4. For any $`F=_\chi A_\chi U_\chi _0^1`$, so that $`\mathrm{\Phi }(F)=_\chi A_\chi Z(\chi )V_\chi _0^2`$, we have the following invariance
$$\mathrm{\Phi }(F),\mathrm{\Phi }(F)_2=\underset{\chi }{}\beta _\chi ^1\left(Z(\chi )^{}A_\chi ^{}A_\chi Z(\chi )\right)=\underset{\chi }{}\beta _\chi ^1\left(A_\chi ^{}A_\chi \right)=F,F_1.$$
From this and recalling the definition of the C\*–norm again in step 5 of the proof of Theorem 3.4 we immediately get the desired isometry property:
$$\mathrm{\Phi }(F)_{}=\underset{\begin{array}{c}X_2_0^2\\ \left|X_2\right|1\end{array}}{sup}|\mathrm{\Phi }(F)X_2|=\underset{\begin{array}{c}X_1_0^1\\ \left|X_1\right|1\end{array}}{sup}|\mathrm{\Phi }(F)X_1|=F_{}.$$
To prove the converse implication assume that $`\mathrm{\Phi }:_1_2`$ specifies the isomorphy of the Hilbert extensions. Use the unitaries $`\{U_\chi \chi 𝒳\}`$ and $`\{V_\chi \chi 𝒳\}`$ in $`_1`$ resp. $`_2`$ to define the unitary
$$Z(\chi ):=\mathrm{\Phi }(U_\chi )V_\chi ^{},\chi 𝒳,$$
that satisfies $`Z(\iota )=\text{}`$. Even more $`Z(\chi )𝒰(𝒵)`$, since for any $`A𝒜`$ we have
$$AZ(\chi )=\mathrm{\Phi }(AU_\chi )V_\chi ^{}=\mathrm{\Phi }\left(U_\chi \beta _\chi ^1(A)\right)V_\chi ^{}=\mathrm{\Phi }(U_\chi )(A^{}V_\chi )^{}=Z(\chi )A.$$
Finally, for $`\chi ,\chi ^{}𝒳`$ we have
$`Z(\chi \chi ^{})`$ $`=`$ $`\mathrm{\Phi }\left(\omega _1(\chi ,\chi ^{})^1U_\chi U_\chi ^{}\right)V_\chi ^{}^{}V_\chi ^{}(\omega _2(\chi ,\chi ^{})^1)^{}`$
$`=`$ $`\omega _1(\chi ,\chi ^{})^1\mathrm{\Phi }(U_\chi )Z(\chi ^{})V_\chi ^{}\omega _2(\chi ,\chi ^{})`$
$`=`$ $`\omega _1(\chi ,\chi ^{})^1\mathrm{\Phi }(U_\chi )\left(\beta _\chi (Z(\chi ^{})^{})V_\chi \right)^{}\omega _2(\chi ,\chi ^{})`$
$`=`$ $`\omega _1(\chi ,\chi ^{})^1Z(\chi )\beta _\chi (Z(\chi ^{}))\omega _2(\chi ,\chi ^{}).`$
Now recalling the definition of the coboundary $`Z`$, the preceding equations imply that $`\omega _1(\chi ,\chi ^{})=Z(\chi ,\chi ^{})\omega _2(\chi ,\chi ^{})`$, $`\chi ,\chi ^{}𝒳`$, and the prove is concluded.
###### Remark 3.9
* Note that the results are independent of the choice of the system $`\beta _𝒳`$ of representatives of $`\mathrm{\Theta }(𝒳)`$. Theorem 3.8 means that there is a bijection between $`H^2(𝒳,𝒰(𝒵),\beta _𝒳)`$ and the set of all $`𝒜`$–module isomorphy classes of Hilbert extensions of $`\{𝒜,\mathrm{\Theta }(𝒳)\}`$ if there is one extension. In other words, the theorem gives an outer characterization of $`H^2(𝒳,𝒰(𝒵),\beta _𝒳)`$ by the set of all $`𝒜`$–module isomorphy classes of Hilbert extensions.
* For a closer analysis of the second cohomology in the special cases were $`\mathrm{\Gamma }\mathrm{}_{\mathrm{}}`$ and $`\mathrm{\Gamma }\mathrm{}_{\mathrm{}}\times \mathrm{}_{\mathrm{}}`$ see . Consider also the abstract results in \[18, Chapter 4\].
## 4 The case of a trivial center
In this case we have $`𝒵=\mathrm{}\text{}`$, thus $`𝒰(𝒵)=\mathrm{𝕋}\text{}`$ and this implies that two automorphisms $`\alpha ,\beta \mathrm{\Gamma }`$ are either unitarily equivalent or otherwise disjoint. The following result is a special case of the famous Doplicher/Roberts theorem (see ) in the present automorphism context.
###### Proposition 4.1
If there is a system of representatives $`ϵ(\alpha ,\beta )`$ of the permutator classes $`\widehat{ϵ}(\alpha ,\beta )`$ which satisfy the equations
$`ϵ(\gamma _1,\gamma _2)ϵ(\gamma _2,\gamma _1)`$ $`=`$ $`\text{},`$
$`ϵ(\iota ,\gamma )=ϵ(\gamma ,\iota )`$ $`=`$ $`\text{},`$
$`\gamma _1(ϵ(\gamma _2,\gamma _3))ϵ(\gamma _1,\gamma _3)`$ $`=`$ $`ϵ(\gamma _1\gamma _2,\gamma _3),`$
$`A\beta _{\chi _1}(B)ϵ(\chi _1,\chi _2)`$ $`=`$ $`ϵ^{}(\chi _1,\chi _2)B\beta _{\chi _2}(A),`$
for all $`A(\beta _{\chi _1},\beta _{\chi _1}^{})`$, $`B(\beta _{\chi _2},\beta _{\chi _2}^{})`$, where $`ϵ^{}`$ belongs to $`\beta _𝒳^{}`$, then there is a generalized 2–cocycle $`\omega _0`$ w.r.t. some system $`\beta _\chi `$ of representatives of the classes $`\chi \mathrm{\Gamma }/\mathrm{int}𝒜`$, with
$$\omega _0(\chi _1,\chi _2)\omega _0(\chi _2,\chi _1)^1=ϵ(\beta _{\chi _1},\beta _{\chi _2}).$$
In this case there is a Hilbert extension $``$ of $`\{𝒜,\mathrm{\Gamma }\}`$.
Conversely, if there is a Hilbert extension $``$ of $`\{𝒜,\mathrm{\Gamma }\}`$, then to each $`\alpha \mathrm{\Gamma }`$ there corresponds a unitary $`V_\alpha _{\chi 𝒳}𝒰(\mathrm{\Pi }_\chi )`$, such that $`\alpha =\text{ad}V_\alpha \text{}𝒜`$ and
$$ϵ(\alpha ,\beta ):=V_\alpha V_\beta V_\alpha ^1V_\beta ^1,$$
is a system of representatives of the permutators $`\widehat{ϵ}(\alpha ,\beta )`$ satisfying the equations above.
###### Remark 4.2
* In the present case the 2-cocycles $`\lambda `$ of the preceding section are $`\mathrm{𝕋}\text{}`$-valued and the relation (6) becomes the usual cocycle equation
$$\lambda (\chi _1,\chi _2)\lambda (\chi _1\chi _2,\chi _3)=\lambda (\chi _2,\chi _3)\lambda (\chi _1,\chi _2\chi _3).$$
* In the particular case where $`𝒜`$ is the inductive limit of a net of von Neumann algebras (which is a standard situation in algebraic quantum field theory, $`𝒜`$ being the so–called quasilocal algebra) it can be shown that there is a lift $`\gamma _𝒳`$ of a given system of representatives $`\beta _𝒳`$, $`\beta _\chi \chi `$ (cf. Definition 2.7), and by Corollary 3.5 we have that $`\omega (\chi _1,\chi _2)=1`$ is an admissible 2–cocycle of the system $`\gamma _𝒳`$. For a detailed construction of the lift see , \[2, Section 3.2\].
## 5 A Hilbert space representation of $`\{,\alpha _𝒢\}`$
Following Sutherland one can introduce a faithful Hilbert space representation of a Hilbert extension $`\{,\alpha _𝒢\}`$ of $`\{𝒜,\mathrm{\Theta }(𝒳)\}.`$
First let $``$ be a Hilbert space and let $`\pi `$ be a faithful representation of $`𝒜`$ on $``$. Form the Hilbert space $`𝒦:=l^2(𝒳,)`$ by completion of $`C_0(𝒳)`$ w.r.t. the norm $`f^2:=_\chi f(\chi )_{}^2`$. Choose a system $`\beta (𝒳)`$ of representatives of $`\mathrm{\Theta }(𝒳)`$ and let $`\omega `$ be a corresponding generalized 2-cocycle such that $`U_{\chi _1}U_{\chi _2}=\omega (\chi _1,\omega _2)U_{\chi _1\chi _2}.`$ Now define a representation $`\mathrm{\Phi }`$ of $`_0`$ on $`𝒦`$ by
$`(\mathrm{\Phi }(A)f)(\chi )`$ $`:=`$ $`\pi (\beta _{\chi ^1}(A))f(\chi ),A𝒜,`$
$`\mathrm{\Phi }(U_{\chi _0})f)(\chi )`$ $`:=`$ $`\pi (\omega (\chi ^1,\chi _0))f(\chi _0^1\chi ),\chi _0𝒳,`$
$`\mathrm{\Phi }(AU_\chi )`$ $`:=`$ $`\mathrm{\Phi }(A)\mathrm{\Phi }(U_\chi ),A𝒜,\chi 𝒳.`$
Note that $`\mathrm{\Phi }(\text{})=\text{}_𝒦`$ and $`\mathrm{\Phi }(A)_𝒦=A.`$ One calculates easily
$`\mathrm{\Phi }(U_{\chi _1})\mathrm{\Phi }(U_{\chi _2})`$ $`=`$ $`\mathrm{\Phi }(\omega (\chi _1,\chi _2))\mathrm{\Phi }(U_{\chi _1\chi _2}),`$
$`\mathrm{\Phi }(U_\chi )\mathrm{\Phi }(A)`$ $`=`$ $`\mathrm{\Phi }(\beta _\chi (A))\mathrm{\Phi }(U_\chi ),`$
$`\mathrm{\Phi }(A^{})=\mathrm{\Phi }(A)^{},\mathrm{\Phi }(U_\chi ^{})`$ $`=`$ $`\mathrm{\Phi }(U_\chi )^{}.`$
Further $`\mathrm{\Phi }(_\chi A_\chi U_\chi )=0`$ implies $`_\chi A_\chi U_\chi =0`$, i.e. $`\mathrm{\Phi }`$ is a \*-isomorphism from $`_0`$ onto $`\mathrm{\Phi }(_0)(𝒦).`$ Recall that
$$\mathrm{\Phi }(F)_𝒦=\text{sup}_{f1}\mathrm{\Phi }(F)f_𝒦.$$
We have
###### Lemma 5.1
The relation
$$\underset{g𝒢}{sup}\mathrm{\Phi }(\alpha _gF)_𝒦<\mathrm{},F_0,$$
(8)
holds.
Proof: With $`F=_\chi A_\chi U_\chi `$ we have
$`\mathrm{\Phi }(F)f^2`$ $`=`$ $`{\displaystyle \underset{y𝒳}{}}{\displaystyle \underset{\chi }{}}\pi (\alpha _{y^1}(A_\chi )\omega (y^1,\chi ))f(y^1\chi )_{}^2`$
$``$ $`{\displaystyle \underset{y𝒳}{}}({\displaystyle \underset{\chi }{}}\pi (\alpha _{y^1}(A_\chi )\omega (y^1,\chi ))f(y^1\chi ))^2`$
$``$ $`{\displaystyle \underset{y𝒳}{}}({\displaystyle \underset{\chi }{}}A_\chi f(y^1\chi ))^2{\displaystyle \underset{y𝒳}{}}({\displaystyle \underset{\chi }{}}A_\chi ^2)({\displaystyle \underset{\chi }{}}f(y^1\chi )^2)`$
$`=`$ $`({\displaystyle \underset{\chi }{}}A_\chi ^2){\displaystyle \underset{\chi }{}}{\displaystyle \underset{y𝒳}{}}f(y^1\chi )^2=N(F)f^2{\displaystyle \underset{\chi }{}}A_\chi ^2,`$
where $`N(F)`$ denotes the number of terms of $`F`$. Hence we obtain
$$\mathrm{\Phi }(F)_𝒦N(F)^{1/2}(\underset{\chi }{}A_\chi ^2)^{1/2}=:C_F.$$
and this implies
$$\mathrm{\Phi }(\alpha _gF)_𝒦C_F,g𝒢,$$
because the number of terms of $`\alpha _gF`$ equals that of $`F`$ and $`\chi (g)A_\chi =A_\chi .`$ This implies the inequality (8).
This result means that
$$\mathrm{\Phi }(F)_{sup}:=\underset{g𝒢}{sup}\mathrm{\Phi }(\alpha _gF)_𝒦$$
is a C\*-norm on $`_0.`$
###### Theorem 5.2
The relation
$$\mathrm{\Phi }(F)_{sup}=F_{},F_0,$$
holds, and in particular $`\mathrm{\Phi }(F)_𝒦F_{}`$, $`F_0`$.
Proof: The norm $`_0F\mathrm{\Phi }(F)_{sup}`$ has the properties $`\mathrm{\Phi }(A)_{sup}=A`$ for all $`A𝒜`$ and $`\mathrm{\Phi }(\alpha _gF)_{sup}=\mathrm{\Phi }(F)_{sup}`$ for all $`g𝒢.`$ However, according to Doplicher/Roberts \[11, p. 105\] there is at most one C\*-norm on $`_0`$ with the mentioned properties.
###### Remark 5.3
If there is a faithful state $`\varphi _0`$ of $`𝒜`$, then Theorem 5.2 can be improved. In this case
$$\mathrm{\Phi }(F)_𝒦=F_{},F_0,$$
holds. This is implied by the fact that in this case Sutherland’s representation $`\mathrm{\Phi }`$ of $`_0`$ on $`𝒦`$ is unitarily equivalent to the so–called regular representation of $`\{,\alpha _𝒢\}`$ (restricted to $`_0`$) given by the (faithful) GNS-representation $`\pi `$ of $`\{,\alpha _𝒢\}`$ on the GNS-Hilbert space $`_\pi `$ w.r.t. the $`𝒢`$-invariant state $`\varphi (F):=\varphi _0(\mathrm{\Pi }_\iota F),F,`$ such that $`\mathrm{\Phi }(F)_𝒦=\pi (F)__\pi `$ for all $`F_0`$, but $`\pi (F)__\pi =F_{}`$ for all $`F`$ (see, for example, \[2, p. 108 ff.\]).
### Acknowledgements
One of us (H.B.) wants to thank Alan Carey for his kind invitation to the University of Adelaide as well as for valuable suggestions of an earlier version of the manuscript. The other author (F.Ll.) expresses his gratitude to Sergio Doplicher for his hospitality at the ‘Dipartamento di Matematica dell’ Università di Roma ‘La Sapienza” in october ’99. The visit was supported by a EU TMR network “Implementation of concept and methods from Non–Commutative Geometry to Operator Algebras and its applications”, contract no. ERB FMRX-CT 96-0073. |
warning/0002/hep-th0002056.html | ar5iv | text | # 1 Introduction
## 1 Introduction
One of the more intriguing outcomes of recent work on the AdS/CFT correspondence has been a renewed effort to understand how the lower-dimensional gauged supergravities arise as Kaluza-Klein sphere reductions from $`D=11`$ or type IIB supergravity. It was long ago demonstrated how the reductions work at the linearised level, but few complete non-linear results existed. A proof of the consistency of the $`S^7`$ reduction from $`D=11`$ was presented, although the Kaluza-Klein Ansatz for the field-strength sector was not fully explicit . It was generally assumed that the other cases, namely the $`S^4`$ reduction of $`D=11`$, and the $`S^5`$ reduction of type IIB, would be consistent too, but until recently no results for these cases had been obtained. In recent work, a fully explicit reduction Ansatz for the $`SO(5)`$-gauged $`N=4`$ $`D=7`$ case has been obtained .<sup>1</sup><sup>1</sup>1The complete bosonic reduction Ansatz for another case, namely the local $`S^4`$ reduction of massive type IIA supergravity to $`SU(2)`$ gauged $`N=2`$ supergravity in $`D=6`$ has also been obtained . Explicit results have also been obtained for various truncations of the full maximal supergravities. These include truncations to the maximal abelian subgroups $`U(1)^4`$, $`U(1)^3`$ and $`U(1)^2`$ in $`D=4`$, 5 and 7 ; the truncation to $`SU(2)`$-gauged $`N=2`$ in $`D=7`$ ; to $`SU(2)\times U(1)`$ gauged $`N=4`$ in $`D=5`$ ; and to $`SO(4)`$ gauged $`N=4`$ in $`D=4`$ . For many purposes, if the fields that participate in the lower-dimensional solutions of interest lie within these truncated subsectors, the truncated reduction is much easier to use, since it is usually much simpler than the full maximal result.<sup>2</sup><sup>2</sup>2We emphasise that in this discussion we are considering only the “remarkable” Kaluza-Klein sphere reductions, where there is no group-theoretic understanding for why the consistency is achievable. In particular, some of the scalar fields parameterise inhomogeneous distortions of the sphere. These contrast with, for example, torus reductions, where the consistency of the truncation to the massless sector is guaranteed by group theory.
Another truncation that allows for relatively simple although still non-trivial sphere reductions is where one retains only the metric and a certain subset of the scalar fields of the lower-dimensional gauged supergravity; one can keep just certain scalars contained in the $`SL(N,R)/SO(N)`$ subset of the full scalar coset manifold, which can be described by a symmetric tensor $`T_{ij}`$. In $`D=4`$, 5 and 7 these subsets correspond to $`N=8`$, 6 and 5 respectively. Specifically, one can consistently truncate to the diagonal scalars, $`T_{ij}=X_i\delta _{ij}`$, where $`_iX_i=1`$. Thus there are 7, 5 and 4 independent scalars in the $`D=4`$, 5 and 7 cases. As shown in , where the reduction Ansätze for these scalar subsectors were presented, one can actually discuss all cases where the lower dimension $`D`$ is related to $`N`$ by
$$N=\frac{4(D2)}{D3},$$
(1)
corresponding to supersymmetric higher-dimensional theories, in a uniform way. The only integer possibilities are $`(D,N)=(4,8)`$, $`(5,6)`$ and $`(7,5)`$, as listed above. (Some proposals for other scalar truncations were presented recently in .) In , extremal AdS domain wall solutions in these dimensions were derived, with the general set of $`(N1)`$ independent charge parameters. By using the reduction Ansätze to oxidise the solutions to the higher dimensions, it was shown how they can be interpreted as continuous distributions of M-branes or D-branes . (Various special cases were obtained also in .)
Certain consistency checks for the reduction Ansätze presented in were conducted there, but a full demonstration of the consistency was not given. Here, we complete the argument by checking all the higher-dimensional equations of motion, and verifying that indeed they are satisfied by the Ansätze of , if and only if the lower-dimensional equations of motion are satisfied.<sup>3</sup><sup>3</sup>3Note that substituting into the higher-dimensional Lagrangian and integrating out the sphere directions could never, per se, yield a proof of consistency. Of course these calculations would be subsumed by complete demonstrations of the consistency of the maximal supergravity reductions in $`D=4`$, 5 and 7. Such a complete proof exists for $`D=7`$ , and implicitly for $`D=4`$ , but not yet for $`D=5`$. Thus the results presented here provide new and independent evidence for the conjectured consistency in all the cases.
## 2 The Scalar Theories, and the Reduction Ansätze
The truncated lower-dimensional gravity plus scalar theory is described by the following Lagrangian in $`D`$ dimensions :
$$e^1_D=R\frac{1}{2}(\stackrel{}{\phi })^2V,$$
(2)
where the potential $`V`$ is given by
$$V=\frac{1}{2}g^2\left((\underset{i=1}{\overset{N}{}}X_i)^22\underset{i=1}{\overset{N}{}}X_i^2\right).$$
(3)
(In $`D=4`$, 5 and 7, we shall have $`N=8`$, 6 and 5 respectively.) The $`N`$ quantities $`X_i`$ are parameterised in terms of $`(N1)`$ independent dilatonic scalars $`\stackrel{}{\phi }`$ as follows:
$$X_i=e^{\frac{1}{2}\stackrel{}{b}_i\stackrel{}{\phi }},$$
(4)
where the $`\stackrel{}{b}_i`$ satisfy
$$\stackrel{}{b}_i\stackrel{}{b}_j=8\delta _{ij}\frac{8}{N},\underset{i}{}\stackrel{}{b}_i=0,(\stackrel{}{u}\stackrel{}{b}_i)\stackrel{}{b}_i=8\stackrel{}{u},$$
(5)
The middle equation here expresses the fact that the $`N`$ quantities $`X_i`$ are subject to the condition
$$\underset{i=1}{\overset{N}{}}X_i=1.$$
(6)
The last equation in (5) allows us to express the dilatons $`\stackrel{}{\phi }`$ in terms of the $`X_i`$:
$$\stackrel{}{\phi }=\frac{1}{4}\underset{i}{}\stackrel{}{b}_i\mathrm{log}X_i.$$
(7)
The equations of motion for the scalar fields, following from (2), are
$$\text{ }\text{ }\stackrel{}{\phi }=\frac{V}{\stackrel{}{\phi }}.$$
(8)
From (4) it follows that $`X_i/\stackrel{}{\phi }=\frac{1}{2}\stackrel{}{b}_iX_i`$, and hence the equations of motion (8) become
$$\text{ }\text{ }\stackrel{}{\phi }=\frac{1}{2}g^2\underset{i}{}\stackrel{}{b}_i\left(X_i\underset{j}{}X_j2X_i^2\right).$$
(9)
Note that we can also write the scalar equations of motion as
$$\text{ }\text{ }\mathrm{log}X_i=2g^2\left(2X_i^2X_i\underset{j}{}X_j\frac{2}{N}\underset{k}{}X_k^2+\frac{1}{N}(\underset{j}{}X_j)^2\right).$$
(10)
The Einstein equation following from (2) is
$$R_{\mu \nu }=\frac{1}{4}X_i^2_\mu X_i_\nu X_i+\frac{1}{D2}Vg_{\mu \nu }.$$
(11)
The Kaluza-Klein sphere reduction Ansätze for obtaining these theories from the higher dimension were presented in , and are as follows:
$`d\widehat{s}^2`$ $`=`$ $`\mathrm{\Delta }^{\frac{2}{D1}}ds_D^2+{\displaystyle \frac{1}{g^2}}\mathrm{\Delta }^{\frac{D3}{D1}}{\displaystyle \underset{i}{}}X_i^1d\mu _i^2,`$
$`\widehat{F}`$ $`=`$ $`g{\displaystyle \underset{i}{}}(2X_i^2\mu _i^2\mathrm{\Delta }X_i)ϵ_{\left(D\right)}+{\displaystyle \frac{1}{2g}}{\displaystyle \underset{i}{}}X_i^1dX_id(\mu _i^2),`$ (12)
where
$$\mathrm{\Delta }=\underset{i}{}X_i\mu _i^2,$$
(13)
and the $`\mu _i`$ are a set of $`N`$ “direction cosines” that satisfy the constraint
$$\underset{i}{}\mu _i^2=1.$$
(14)
In (12), $`ϵ_{\left(D\right)}`$ denotes the volume form of the $`D`$-dimensional metric $`ds_D^2`$. Note that if all the scalars $`X_i`$ are trivial ($`X_i=1`$), the internal part of the metric becomes $`_id\mu _i^2`$, which is the metric on the unit $`(N1)`$-sphere.
The $`D`$-form field strength $`\widehat{F}`$ in (12) is the 4-form of eleven-dimensional supergravity for the case $`D=4`$, the Hodge dual of this 4-form for the case $`D=7`$, and it is the self-dual 5-form of the type IIB theory when $`D=5`$. Note that in each case, given the nature of the Ansatz, the relevant Bianchi identity and field equation for $`\widehat{F}`$ are simply
$$d\widehat{F}=0,d\widehat{}\widehat{F}=0.$$
(15)
## 3 The Consistency of the Reduction
It was shown in that the $`D`$-form field-strength Ansatz in (12) satisfies the Bianchi identity $`d\widehat{F}=0`$, provided that the scalar fields $`X_i`$ satisfy precisely the lower-dimensional equations of motion (10). This calculation is a straightforward one, and we shall not repeat it here. It is harder to show that $`\widehat{F}`$ satisfies the field equation $`d\widehat{}\widehat{F}=0`$, because this involves taking a Hodge dual of the field strength $`\widehat{F}`$. This is what we shall now address.
### 3.1 The Field Equation for $`\widehat{F}`$
The complication here is that the $`(N1)`$-sphere is being coordinatised by $`N`$ quantities $`\mu _i`$ subject to the constraint (14). It seems that the best way to proceed is to eliminate one of the $`\mu _i`$ in favour of the others, using (14). To that end, we split the $`\mu _i`$ as $`\mu _i=(\mu _\alpha ,\mu _0)`$, and solve for $`\mu _0`$ in terms of the $`\mu _\alpha `$.
If we consider first the metric
$$ds^2=\underset{i}{}X_i^1d\mu _i^2,$$
(16)
then in terms of the $`\mu _\alpha `$ we can write it as $`ds^2=g_{\alpha \beta }d\mu _\alpha d\mu _\beta `$, where
$$g_{\alpha \beta }=X_\alpha ^1\delta _{\alpha \beta }+\frac{1}{X_0\mu _0^2}\mu _\alpha \mu _\beta .$$
(17)
(We have, of course, used the identity
$$d\mu _0=\frac{\mu _\alpha }{\mu _0}d\mu _\alpha ,$$
(18)
which follows from (14).)
It is straightforward to invert the metric $`g_{\alpha \beta }`$ given in (17). The result is
$$g^{\alpha \beta }=X_\alpha \delta _{\alpha \beta }\mathrm{\Delta }^1\mu _\alpha \mu _\beta X_\alpha X_\beta .$$
(19)
It is also easy to establish that
$$det(g_{\alpha \beta })=\frac{\mathrm{\Delta }}{\mu _0^2}.$$
(20)
Note that it follows from the metric Ansatz in (12) that the determinant of the higher-dimensional metric $`d\widehat{s}^2`$ is given by
$$det(\widehat{g})=\frac{\mathrm{\Delta }^{\frac{4}{D1}}}{g^{2N2}\mu _0^2}det(g_D),$$
(21)
where $`g_D`$ denotes the $`D`$-dimensional spacetime metric $`ds_D^2`$ and $`g`$ in the denominator is just the gauge coupling constant (not to be confused with the determinant of the higher-dimensional metric $`\widehat{g}`$ or the one for the lower dimension, $`g_D`$.)
Now let us look at the field-strength Ansatz. We shall use the convention that $`\epsilon _{M_1\mathrm{}M_D}`$ always means the tensor density, which is the pure numbers $`\pm 1,0`$. So the Ansatz for $`\widehat{F}`$ in (12) is
$`\widehat{F}_{\nu _1\mathrm{}\nu _D}`$ $`=`$ $`gU\sqrt{g_D}\epsilon _{\nu _1\mathrm{}\nu _D},`$
$`\widehat{F}_{\nu _1\mathrm{}\nu _{D1}\alpha }`$ $`=`$ $`{\displaystyle \frac{1}{g}}\sqrt{g_D}\epsilon _{\nu _1\mathrm{}\nu _{D1}\rho }g_D^{\rho \sigma }(X_\alpha ^1_\sigma X_\alpha X_0^1_\sigma X_0)\mu _\alpha ,`$ (22)
where
$$U\underset{i}{}(2X_i^2\mu _i^2\mathrm{\Delta }X_i),$$
(23)
and $`g_D`$ denotes the $`D`$-dimensional spacetime metric $`ds_D^2`$.
We can now calculate the upper-index components of $`\widehat{F}`$. In fact, what we really need is these components multiplied by $`\sqrt{\widehat{g}}`$. From the results above we find
$`\sqrt{\widehat{g}}\widehat{F}^{\nu _1\mathrm{}\nu _D}`$ $`=`$ $`{\displaystyle \frac{U}{g^{N2}\mu _0\mathrm{\Delta }^2}}\epsilon ^{\nu _1\mathrm{}\nu _D},`$
$`\sqrt{\widehat{g}}\widehat{F}^{\nu _1\mathrm{}\nu _{D1}\alpha }`$ $`=`$ $`{\displaystyle \frac{1}{g^{N2}\mu _0}}\epsilon ^{\nu _1\mathrm{}\nu _{D1}\sigma }_\sigma \left({\displaystyle \frac{X_\alpha \mu _\alpha }{\mathrm{\Delta }}}\right).`$ (24)
($`\epsilon ^{M_1\mathrm{}M_D}`$ is the tensor density that takes the values 0, $`\pm 1`$, and is numerically equal to $`\epsilon _{M_1\mathrm{}M_D}`$.)
One can directly verify from these expressions that the field equation is satisfied, namely that
$$_M\left(\sqrt{\widehat{g}}\widehat{F}^{N_1\mathrm{}N_{D1}M}\right)=0.$$
(25)
However, it is more elegant to do this by using (24) first to construct the Hodge dual of $`\widehat{F}`$ itself. To do this, we make the following definitions:
$`P`$ $``$ $`{\displaystyle \frac{1}{n!}}\epsilon _{\alpha _1\mathrm{}\alpha _n}d\mu _{\alpha _1}\mathrm{}d\mu _{\alpha _n},`$
$`Q_\alpha `$ $``$ $`{\displaystyle \frac{1}{(n1)!}}\epsilon _{\alpha \beta _1\mathrm{}\beta _{n1}}d\mu _{\beta _1}\mathrm{}d\mu _{b_{n1}},`$
$`W`$ $``$ $`{\displaystyle \frac{1}{n!}}\epsilon _{ij_1\mathrm{}j_n}\mu _id\mu _{j_1}\mathrm{}d\mu _{j_n},`$
$`Z_i`$ $``$ $`{\displaystyle \frac{1}{(n1)!}}\epsilon _{ijk_1\mathrm{}k_{n1}}\mu _jd\mu _{k_1}\mathrm{}d\mu _{k_{n1}},`$ (26)
where $`n=N1`$. Note that what we have done here is to define $`P`$ and $`Q_\alpha `$ with respect to the reduced set of $`n=N1`$ coordinates $`\mu _\alpha `$, while $`W`$ and $`Z_i`$ are defined with respect to the full set of $`N`$ coordinates $`\mu _i`$. (Some analogous formulae and manipulations are presented also in .)
Now, we can establish the following:
$`W`$ $`=`$ $`{\displaystyle \frac{1}{\mu _0}}P,`$
$`Z_0`$ $`=`$ $`\mu _\beta Q_\beta ,Z_\alpha ={\displaystyle \frac{1}{\mu _0}}(Q_\alpha +\mu _\alpha \mu _\beta Q_\beta ),`$
$`d\mu _\alpha Q_\beta `$ $`=`$ $`P\delta _{\alpha \beta },`$
$`d\mu _iZ_j`$ $`=`$ $`(\delta _{ij}\mu _i\mu _j)W,`$
$`dQ_\alpha `$ $`=`$ $`0,dW=0,dZ_i=n\mu _iW.`$ (27)
From (24), it is evident that we have
$$\widehat{}\widehat{F}_{\alpha _1\mathrm{}\alpha _n}=\frac{U}{g^{N2}\mu _0\mathrm{\Delta }^2}\epsilon _{\alpha _1\mathrm{}\alpha _n},\widehat{}\widehat{F}_{\alpha _1\mathrm{}\alpha _{n1}\nu }=\frac{1}{g^{N2}\mu _0}\epsilon _{\alpha _1\mathrm{}\alpha _{n1}\beta }_\nu \left(\frac{X_\beta \mu _\beta }{\mathrm{\Delta }}\right).$$
(28)
Note that here, and in many other formulae, we are using a “generalised Einstein summation convention,” in which any dummy index that appears two or more times in an expression is understood to be summed over. It will always be clear from context whether an index is a dummy or not.
After some algebra, we can show from the above definitions and properties that this can be written as
$$\widehat{}\widehat{F}=\frac{1}{g^{N2}}\left(\frac{U}{\mathrm{\Delta }^2}W+_\nu \left(\frac{X_i\mu _i}{\mathrm{\Delta }}\right)dx^\nu Z_i\right).$$
(29)
To check that the equation of motion $`d\widehat{}\widehat{F}=0`$ is satisfied, we just have to make use of the various lemmata established above. Thus we have
$`g^{N2}d\widehat{}\widehat{F}`$ $`=`$ $`_\nu \left({\displaystyle \frac{U}{\mathrm{\Delta }^2}}\right)dx^\nu W_\nu \left({\displaystyle \frac{X_i\mu _i}{\mathrm{\Delta }}}\right)dx^\nu dZ_i_{\mu _j}_\nu \left({\displaystyle \frac{X_i\mu _i}{\mathrm{\Delta }}}\right)dx^\nu d\mu _jZ_i,`$ (30)
$`=`$ $`_\nu \left({\displaystyle \frac{U}{\mathrm{\Delta }^2}}\right)dx^\nu Wn\mu _i_\nu \left({\displaystyle \frac{X_i\mu _i}{\mathrm{\Delta }}}\right)dx^\nu W`$
$`+_{\mu _j}_\nu \left({\displaystyle \frac{X_i\mu _i}{\mathrm{\Delta }}}\right)dx^\nu W(\delta _{ij}\mu _i\mu _j),`$
$`=`$ $`_\nu \left({\displaystyle \frac{U}{\mathrm{\Delta }^2}}\right)dx^\nu W+_\nu \left(\delta _{ij}{\displaystyle \frac{X_i}{\mathrm{\Delta }}}{\displaystyle \frac{2X_iX_j\mu _i\mu _j}{\mathrm{\Delta }^2}}\right)dx^\nu W(\delta _{ij}\mu _i\mu _j),`$
$`=`$ $`_\nu \left({\displaystyle \frac{U}{\mathrm{\Delta }^2}}\right)dx^\nu W_\nu \left({\displaystyle \frac{U}{\mathrm{\Delta }^2}}\right)dx^\nu W=0.`$
Note that in various steps above, we have made use of the fact that the $`\mu _i`$ can be taken freely inside the $`_\nu `$ derivative, and that therefore, for instance, a term like $`\mu _i_\nu (X_i\mu _i/\mathrm{\Delta })`$ is equal to $`_\nu (X_i\mu _i^2/\mathrm{\Delta })`$, which is therefore zero since $`X_i\mu _i^2=\mathrm{\Delta }`$. This completes the checking of the consistency of the higher-dimensional field equation for $`\widehat{F}`$.
### 3.2 The Einstein Equation
#### 3.2.1 Calculation of the Ricci Tensor
To check the various components of the higher-dimensional Einstein equation, we first calculate the curvature tensor for the metric Ansatz. From now on, since no generality is lost, we set the gauge coupling $`g`$ equal to 1 for simplicity. The metric can be written as
$$d\widehat{s}^2=\mathrm{\Delta }^ads_D^2+\mathrm{\Delta }^b\underset{i}{}X_i^1d\mu _i^2,$$
(31)
where
$$a=\frac{2}{D1},b=\frac{D3}{D1}.$$
(32)
From this, we find that the affine connection $`\widehat{\mathrm{\Gamma }}^M{}_{NP}{}^{}=\frac{1}{2}\widehat{g}^{MQ}(_Ng_{QP}+_Pg_{QN}_Qg_{NP})`$ is given by
$`\widehat{\mathrm{\Gamma }}^\mu _{\nu \rho }`$ $`=`$ $`\mathrm{\Gamma }^\mu {}_{\nu \rho }{}^{}+\frac{1}{2}a\mathrm{\Delta }^1(\delta _\rho ^\mu _\nu \mathrm{\Delta }+\delta _\nu ^\mu _\rho \mathrm{\Delta }g_{\nu \rho }^\mu \mathrm{\Delta }),`$
$`\widehat{\mathrm{\Gamma }}^\mu _{\nu \alpha }`$ $`=`$ $`\frac{1}{2}a\mathrm{\Delta }^1\delta _\nu ^\mu _\alpha \mathrm{\Delta },`$
$`\widehat{\mathrm{\Gamma }}^\alpha _{\mu \nu }`$ $`=`$ $`\frac{1}{2}ag_{\mu \nu }^\alpha \mathrm{\Delta },`$
$`\widehat{\mathrm{\Gamma }}^\alpha _{\beta \mu }`$ $`=`$ $`\frac{1}{2}b\mathrm{\Delta }^1\delta _\beta ^\alpha _\mu \mathrm{\Delta }+\frac{1}{2}g^{\alpha \gamma }_\mu g_{\beta \gamma },`$
$`\widehat{\mathrm{\Gamma }}^\mu _{\alpha \beta }`$ $`=`$ $`\frac{1}{2}bg_{\alpha \beta }\mathrm{\Delta }^2^\mu \mathrm{\Delta }\frac{1}{2}\mathrm{\Delta }^1^\mu g_{\alpha \beta },`$
$`\widehat{\mathrm{\Gamma }}^\alpha _{\beta \gamma }`$ $`=`$ $`\mathrm{\Gamma }^\alpha {}_{\beta \gamma }{}^{}\frac{1}{2}b\mathrm{\Delta }^1(\delta _\gamma ^\alpha _\beta \mathrm{\Delta }+\delta _\beta ^\alpha _\gamma \mathrm{\Delta }g_{\beta \gamma }^\alpha \mathrm{\Delta }),`$ (33)
where
$$\mathrm{\Gamma }^\alpha {}_{\beta \gamma }{}^{}\frac{1}{2}g^{\alpha \delta }(_\beta g_{\delta \gamma }+_\gamma g_{\delta \beta }_\delta g_{\beta \gamma })=\mathrm{\Delta }^1X_\alpha \mu _\alpha (\delta _{\beta \gamma }+\widehat{\mu }_\beta \widehat{\mu }_\gamma ).$$
(34)
Note that $`_\alpha `$ means $`/\mu _\alpha `$, and that $`^\alpha g^{\alpha \beta }_\beta `$.
We calculate the curvature using the expressions
$`\widehat{R}^M_{NPQ}`$ $`=`$ $`_P\widehat{\mathrm{\Gamma }}^M{}_{NQ}{}^{}_Q\widehat{\mathrm{\Gamma }}^M{}_{NP}{}^{}+\widehat{\mathrm{\Gamma }}^M{}_{PR}{}^{}\widehat{\mathrm{\Gamma }}_{}^{R}{}_{QN}{}^{}\widehat{\mathrm{\Gamma }}^M{}_{QR}{}^{}\widehat{\mathrm{\Gamma }}_{}^{R}{}_{PN}{}^{},`$
$`\widehat{R}_{NQ}`$ $``$ $`\widehat{R}^M{}_{NMQ}{}^{}=_M\widehat{\mathrm{\Gamma }}^M{}_{NQ}{}^{}_Q\widehat{\mathrm{\Gamma }}^M{}_{NM}{}^{}+\widehat{\mathrm{\Gamma }}^M{}_{MR}{}^{}\widehat{\mathrm{\Gamma }}_{}^{R}{}_{QN}{}^{}\widehat{\mathrm{\Gamma }}^M{}_{QR}{}^{}\widehat{\mathrm{\Gamma }}_{}^{R}{}_{MN}{}^{}.`$ (35)
After some calculation, we find that
$`\widehat{R}_{\mu \nu }`$ $`=`$ $`R_{\mu \nu }\frac{1}{4}X_i^2_\mu X_i_\nu X_i+\frac{1}{2}\mathrm{\Delta }^1X_i^1\mu _i^2_\mu X_i_\nu X_i\frac{1}{2}\mathrm{\Delta }^2_\mu \mathrm{\Delta }_\nu \mathrm{\Delta }`$ (36)
$`+\frac{1}{2}a(\mathrm{\Delta }^2_\lambda \mathrm{\Delta }^\lambda \mathrm{\Delta }\mathrm{\Delta }^1\text{ }\text{ }\mathrm{\Delta })g_{\mu \nu }`$
$`a\left[{\displaystyle \underset{i}{}}X_i^2\mathrm{\Delta }^1X_i^2\mu _i^2{\displaystyle \underset{j}{}}X_j2\mathrm{\Delta }^1X_i^3\mu _i^2+2\mathrm{\Delta }^2(X_i^2\mu _i^2)^2\right]g_{\mu \nu },`$
$`\widehat{R}_{\alpha \beta }`$ $`=`$ $`R_{\alpha \beta }+\frac{1}{2}bg_{\alpha \beta }\mathrm{\Delta }^2\text{ }\text{ }\mathrm{\Delta }\frac{1}{2}bg_{\alpha \beta }\mathrm{\Delta }^3_\lambda \mathrm{\Delta }^\lambda \mathrm{\Delta }\frac{1}{2}\mathrm{\Delta }^1\text{ }\text{ }g_{\alpha \beta }`$ (37)
$`+\frac{1}{2}\mathrm{\Delta }^1g^{\gamma \delta }_\lambda g_{\alpha \gamma }^\lambda g_{\beta \delta }\frac{1}{4}\mathrm{\Delta }^2_\alpha \mathrm{\Delta }_\beta \mathrm{\Delta }\frac{1}{2}\mathrm{\Delta }^1_\alpha _\beta \mathrm{\Delta }`$
$`\frac{1}{4}bg_{\alpha \beta }\mathrm{\Delta }^2_\gamma \mathrm{\Delta }^\gamma \mathrm{\Delta }+\frac{1}{2}bg_{\alpha \beta }\mathrm{\Delta }^1_\gamma ^\gamma \mathrm{\Delta },`$
$`\widehat{R}_{\alpha \mu }`$ $`=`$ $`\frac{1}{2}\mathrm{\Delta }^2U(X_\alpha ^1_\mu X_\alpha X_0^1_\mu X_0)\mu _\alpha .`$ (38)
Note that in these expressions $`\text{ }\text{ }`$ means the Laplacian in the lower-dimensional spacetime, $`_\alpha `$ denotes the covariant derivative with respect to the internal metric $`g_{\alpha \beta }`$, with its affine connection $`\mathrm{\Gamma }^\gamma _{\alpha \beta }`$, and $`R_{\alpha \beta }`$ is the Ricci tensor calculated in this connection.
Some useful lemmata which we used are
$`^\alpha \mathrm{\Delta }_\alpha \mathrm{\Delta }`$ $`=`$ $`4X_i^3\mu _i^24\mathrm{\Delta }^1(X_i^2\mu _i^2)^2,`$
$`\mathrm{\Gamma }^\alpha _{\alpha \beta }`$ $`=`$ $`\frac{1}{2}\mathrm{\Delta }^1_\beta \mathrm{\Delta }+{\displaystyle \frac{1}{\mu _0^2}}\mu _\beta ,`$
$`_\alpha ^\alpha \mathrm{\Delta }`$ $`=`$ $`2{\displaystyle \underset{i}{}}X_i^22\mathrm{\Delta }^1X_i^2\mu _i^2{\displaystyle \underset{j}{}}X_j+4\mathrm{\Delta }^2(X_i^2\mu _i^2)^2`$
$`4\mathrm{\Delta }^1X_i^3\mu _i^2+\frac{1}{2}\mathrm{\Delta }^1^a\mathrm{\Delta }_\alpha \mathrm{\Delta },`$
$`R_{\alpha \beta }`$ $`=`$ $`\mathrm{\Delta }^1\overline{g}_{\alpha \beta }{\displaystyle \underset{\gamma }{}}X_\gamma \mathrm{\Delta }^2(X_i^2\mu _i^2)\overline{g}_{\alpha \beta }`$
$`+\mathrm{\Delta }^2(X_\alpha X_0)(X_\beta X_0)\mu _\alpha \mu _\beta \mathrm{\Delta }^1(X_\alpha X_0)\delta _{\alpha \beta },`$
$`\text{ }\text{ }g_{\alpha \beta }`$ $`=`$ $`X_\alpha ^3_\lambda X_\alpha ^\lambda X^\alpha \delta _{\alpha \beta }+X_0^3_\lambda X_0^\lambda X_0\widehat{\mu }_\alpha \widehat{\mu }_\beta `$
$`4(X_\alpha \delta _{\alpha \beta }+X_0\widehat{\mu }_\alpha \widehat{\mu }_\beta )+2\overline{g}_{\alpha \beta }{\displaystyle \underset{j}{}}X_j+{\displaystyle \frac{4}{N}}Vg_{\alpha \beta },`$
$`g^{\gamma \delta }_\lambda g_{\alpha \gamma }^\lambda g_{\beta \delta }`$ $`=`$ $`X_\alpha ^3_\lambda X_\alpha ^\lambda X_\alpha \delta _{\alpha \beta }+X_0^3_\lambda X_0^\lambda X_0\widehat{\mu }_\alpha \widehat{\mu }_\beta `$
$`\mathrm{\Delta }^1(X_\alpha ^1_\lambda X_\alpha X_0^1_\lambda X_0)(X_\beta ^1^\lambda X_\beta X_0^1^\lambda X_0)\mu _\alpha \mu _\beta .`$
The quantities $`\widehat{\mu }_\alpha `$ are defined by $`\widehat{\mu }_\alpha \mu _\alpha /\mu _0`$, and the metric $`\overline{g}_{\alpha \beta }`$ is defined by
$$\overline{g}_{\alpha \beta }\delta _{\alpha \beta }+\widehat{\mu }_\alpha \widehat{\mu }_\beta .$$
(40)
It is evident from (17) that $`\overline{g}_{\alpha \beta }`$ is the metric on the unit round $`(N1)`$-sphere, corresponding to setting all the $`X_i=1`$.
#### 3.2.2 The Consistency of the Einstein Equation
With the results for the Ricci tensor from the previous section, we can now verify that all components of the higher-dimensional Einstein equation are indeed consistently satisfied.
The higher-dimensional Einstein equation is
$$\widehat{R}_{MN}=\widehat{S}_{MN},$$
(41)
where
$$\widehat{S}_{MN}=\frac{1}{2(D1)!}\left[\widehat{F}_{MN}^2\frac{D3}{D(D1)}\widehat{F}^2\widehat{g}_{MN}\right].$$
(42)
The non-zero components of $`\widehat{F}_{M_1\mathrm{}M_D}`$ are given in (22). After some algebra, we find that
$`\widehat{F}^2`$ $`=`$ $`D!\mathrm{\Delta }^{Da}(U^2+\mathrm{\Delta }X_i^1\mu _i^2_\lambda X_i^\lambda X_i_\lambda \mathrm{\Delta }^\lambda \mathrm{\Delta }),`$
$`\widehat{F}_{\mu \nu }^2`$ $`=`$ $`(D1)!\mathrm{\Delta }^2[\mathrm{\Delta }X_i^1\mu _i^2_\mu X_i_\nu X_i_\mu \mathrm{\Delta }_\nu \mathrm{\Delta }`$ (43)
$`(\mathrm{\Delta }X_i^1\mu _i^2_\lambda X_i^\lambda X_i_\lambda \mathrm{\Delta }^\lambda \mathrm{\Delta })g_{\mu \nu }U^2g_{\mu \nu }],`$
$`\widehat{F}_{\alpha \beta }^2`$ $`=`$ $`(D1)!\mathrm{\Delta }^2(X_\alpha ^1_\lambda X_\alpha X_0^1_\lambda X_0)(X_\beta ^1^\lambda X_\beta X_0^1^\lambda X_0)\mu _\alpha \mu _\beta ,`$
where, as usual, $`U`$ is given by
$$U=2X_i^2\mu _i^2\mathrm{\Delta }\underset{i}{}X_i.$$
(44)
Thus we find that $`\widehat{S}_{MN}`$ is given by
$`\widehat{S}_{\mu \nu }`$ $`=`$ $`\frac{1}{2}\mathrm{\Delta }^1X_i^1\mu _i^2_\mu X_i_\nu X_i\frac{1}{2}\mathrm{\Delta }^2_\mu \mathrm{\Delta }_\nu \mathrm{\Delta }`$ (45)
$`{\displaystyle \frac{1}{D1}}\mathrm{\Delta }^2(U^2_\lambda \mathrm{\Delta }^\lambda \mathrm{\Delta }+\mathrm{\Delta }X_i^1\mu _i^2_\lambda X_i^\lambda X_i)g_{\mu \nu },`$
$`\widehat{S}_{\alpha \beta }`$ $`=`$ $`\frac{1}{2}b\mathrm{\Delta }^3U^2g_{\alpha \beta }+\frac{1}{2}b\mathrm{\Delta }^2g_{\alpha \beta }X_i^1\mu _i^2_\lambda X_i^\lambda X_i\frac{1}{2}b\mathrm{\Delta }^3_\lambda \mathrm{\Delta }^\lambda \mathrm{\Delta }g_{\alpha \beta }`$ (46)
$`\frac{1}{2}\mathrm{\Delta }^2(X_\alpha ^1_\lambda X_\alpha X_0^1_\lambda X_0)(X_\beta ^1^\lambda X_\beta X_0^1^\lambda X_0)\mu _\alpha \mu _\beta ,`$
$`\widehat{S}_{\alpha \mu }`$ $`=`$ $`\frac{1}{2}\mathrm{\Delta }^2U(X_\alpha ^1_\mu X_\alpha X_0^1_\mu X_0)\mu _\alpha .`$ (47)
To verify that the components $`\widehat{R}_{\mu \nu }=\widehat{S}_{\mu \nu }`$ of the higher-dimensional Einstein equation indeed imply the lower-dimensional Einstein equation (11), we simply need to substitute the above results into (41). It is also necessary to use the scalar equations of motion in (10), from which we can deduce that
$$\text{ }\text{ }\mathrm{\Delta }=X_i^1\mu _i^2_\lambda X_i^\lambda X_i+4X_i^3\mu _i^22X_i^2\mu _i^2\underset{j}{}X_j\frac{4}{N}\mathrm{\Delta }V.$$
(48)
Putting all the results together, we find that indeed all the $`\mu _i`$ dependence cancels out in the $`\widehat{R}_{\mu \nu }=\widehat{S}_{\mu \nu }`$ equation, and we correctly reproduce the lower-dimensional Einstein equation in (11).
After some algebra, using the lemmata given previously, we find that the components $`\widehat{R}_{\alpha \beta }`$ of the Ricci tensor of the higher-dimensional metric are simply given by
$`\widehat{R}_{\alpha \beta }`$ $`=`$ $`\frac{1}{2}b\mathrm{\Delta }^3U^2g_{\alpha \beta }+\frac{1}{2}\mathrm{\Delta }^2g_{\alpha \beta }X_i^1\mu _i^2_\lambda X_i^\lambda X_i\frac{1}{2}b\mathrm{\Delta }^3_\lambda \mathrm{\Delta }^\lambda \mathrm{\Delta }`$ (49)
$`\frac{1}{2}\mathrm{\Delta }^1(X_\alpha ^1_\lambda X_\alpha X_0^1_\lambda X_0)(X_\beta ^1^\lambda X_\beta X_0^1^\lambda X_0)\mu _\alpha \mu _\beta .`$
Note that we have made use of the equations of motion for the $`X_i`$ fields in simplifying this expression. It is now straightforward to see that this is exactly equal to the expression for $`\widehat{S}_{\alpha \beta }`$ obtained in (46). Finally, the components $`\widehat{S}_{\alpha \mu }`$ given in (47) agree precisely with the corresponding components $`\widehat{R}_{\alpha \mu }`$ found in (38). Thus the consistency of the reduction Ansatz is completely verified.
## 4 Scalar Potentials in $`D=3`$
In the previous sections, we proved the consistency of the embedding of the diagonal symmetric potentials in the relevant higher dimensions. The number of scalars $`N`$ and the (lower) dimension $`D`$ are related by (1). As was shown in , the various $`D`$-dimensional multi-charge extremal AdS domain walls supported by these scalars can be oxidised back to solutions of eleven-dimensional supergravity ($`D=4`$ and $`D=7`$) or type IIB supergravity $`D=5`$). These higher-dimensional solutions correspond to ellipsoidal continuous distributions of M5-branes, M2-branes and D3-branes respectively .
For general values of $`D`$ the relation (1) would imply a non-integral value for $`N`$, and no consistent embedding exists. The relation becomes singular for the case $`D=3`$. Thus contrary to what one might have hoped, the pattern of consistent embeddings does not seem to extend to an $`S^3`$ reduction from $`D=6`$ to $`D=3`$. Indeed, it is straightforward to show that the ellipsoidal continuous distributions of dyonic strings that exist in $`D=6`$ do not lend themselves to consistent reductions to $`D=3`$.
In this section, we discuss an alternative reduction to a gauged $`D=3`$ supergravity, in which there is a massive scalar field. The three-dimensional bosonic Lagrangian is given by
$$e^1_3=R\frac{1}{2}(\varphi )^2\frac{1}{2}g^2\left(\frac{1}{a_1^2}e^{a_1\varphi }\frac{1}{a_1a_2}e^{a_2\varphi }\right),$$
(50)
where $`a_1^2=4/k+4`$ and $`a_2=4/a_1`$. The integer $`k`$ can take the values 1, 2, or 3. The values $`k=2`$ and $`k=3`$ correspond to the $`S^3`$ reduction of $`D=6`$ simple (chiral) supergravity and the $`S^2`$ reduction of $`D=5`$ simple supergravity respectively, and $`\varphi `$ is the associated massive breathing mode .
The case of $`k=1`$ corresponds to the $`S^1`$ Scherk-Schwarz reduction of the Freedman-Schwarz model. To show this, we begin from the Lagrangian for the gravity plus scalar sector of the $`D=4`$ Freedman-Schwarz model , which can be obtained as a singular limit of the $`N=4`$, $`D=4`$, $`SO(4)`$ gauged supergravity :
$$\widehat{e}^1_4=\widehat{R}\frac{1}{2}(\widehat{\varphi })^2\frac{1}{2}(\chi )^2e^{2\widehat{\varphi }}+\frac{1}{2}g^2e^{\widehat{\varphi }}.$$
(51)
Dimensionally reducing this theory on a coordinate $`z`$, where the axion $`\chi `$ is allowed to take the generalised Scherk-Schwarz form $`\chi =mz`$, we obtain the three-dimensional scalar Lagrangian
$$e^1_3=R\frac{1}{2}(\widehat{\varphi })^2\frac{1}{2}(\phi )^2\frac{1}{2}m^2e^{2(\widehat{\varphi }+\phi )}+\frac{1}{2}g^2e^{\widehat{\varphi }+\phi }.$$
(52)
Since the original dilaton $`\widehat{\varphi }`$ and the dilaton $`\phi `$ coming from the dimensional reduction occur everywhere in the same combination, we see that it is consistent to truncate out the combination $`\widehat{\varphi }\phi `$. Making the redefinition $`\varphi (\widehat{\varphi }+\phi )/\sqrt{2}`$, the Lagrangian (52) reduces to (50) with $`k=1`$. The three Lagrangians in (50) all give rise to supersymmetric domain-wall solutions in $`D=3`$ .
## 5 Conclusion
In this paper, we have provided a complete proof of the consistency of the Kaluza-Klein reduction Ansätze that were presented in , which describe the embedding of certain $`N`$-scalar truncations of the maximal gauged supergravities in $`D=4`$, 7 and 5, via spherical reductions on $`S^7`$, $`S^4`$ and $`S^5`$ respectively. The $`N`$ scalars, with $`N=8`$, 5 and 6 respectively, correspond to the diagonal elements in the $`SL(N,R)/SO(N)`$ submanifolds of the full scalar manifolds in the corresponding maximal supergravities. (Actually, there are really only $`N1`$ independent scalars in these truncations, on account of a unit-determinant condition on the scalars in the coset.) Our proof included a complete verification of the consistency of the reduction of the higher-dimensional Einstein equation, which is usually the most calculationally difficult part of the procedure.
For $`D=7`$, our results are consistent with the full Kaluza-Klein $`S^4`$ reduction that was recently obtained explicitly in . For $`D=4`$, they are compatible with the implicit proof of the consistency of the complete $`S^7`$ reduction, presented in . Furthermore, our results provide a complete proof of the validity of the explicit expressions presented in for the Ansätze for the eleven-dimensional fields, which, especially in the case of the 4-form field strength, are not straightforward to extract from the results presented in .
Finally, in $`D=5`$ our results provide further evidence for the conjectured consistency of the $`S^5`$ reduction of type IIB supergravity to give maximal $`SO(6)`$ gauged supergravity in $`D=5`$.
We also considered the special case of scalar theories in $`D=3`$ that arise from dimensional reduction. This dimension lies outside the set of cases covered by the previous discussion, on account of a degeneration in the formula (1) relating the dimension to the number of scalar fields. Instead, we described the set of three theories (50) arising as the scalar sectors of sphere reductions from $`D=6`$, $`D=5`$ and $`D=4`$. In the case of $`D=4`$, we showed how the single-scalar Lagrangian (50) arises from a Scherk-Schwarz $`S^1`$ reduction of the $`D=4`$ Freedman-Schwarz model, accompanied by a further consistent truncation of one combination of the two resulting dilatonic scalar fields.
## Acknowledgements
We are grateful to Jim Liu and Tuan Tran for helpful discussions. |
warning/0002/hep-lat0002010.html | ar5iv | text | # One-loop renormalization of fermionic currents with the overlap-Dirac operator
## I Introduction
Recent developments in lattice QCD have shown that chiral symmetry can be realized on the lattice without fermion doubling (for recent reviews see e.g. Refs.), circumventing the Nielsen-Ninomiya theorem . This has been achieved by introducing an overlap-Dirac operator derived from the overlap formulation of chiral fermions . The simplest such example, for a massless fermion, is given by the Neuberger-Dirac operator
$`D_\mathrm{N}`$ $`=`$ $`{\displaystyle \frac{1}{a}}\rho \left[1+X(X^{}X)^{1/2}\right],`$ (1)
$`X`$ $`=`$ $`D_\mathrm{W}{\displaystyle \frac{1}{a}}\rho ,`$ (2)
where $`D_\mathrm{W}`$ is the Wilson-Dirac operator (with the Wilson parameter $`r`$ set to its standard value, $`r=1`$)
$$D_\mathrm{W}=\frac{1}{2}\left[\gamma _\mu \left(_\mu ^{}+_\mu \right)a_\mu ^{}_\mu \right],$$
(3)
$$_\mu \psi (x)=\frac{1}{a}\left[U(x,\mu )\psi (x+a\widehat{\mu })\psi (x)\right],$$
(4)
$`a`$ is the lattice spacing and $`\rho `$ is a real parameter subject to the constraint $`0<\rho <2`$. Nonperturbatively one expects $`m_c<\rho <2`$, where $`m_c<0`$ is the critical mass associated with the Wilson-Dirac operator.
$`D_\mathrm{N}`$ satisfies the Ginsparg-Wilson relation
$$\gamma _5D+D\gamma _5=aD\gamma _5D,$$
(5)
which protects the quark masses from additive renormalization . Lattice gauge theories with Ginsparg-Wilson fermions have been proved to be renormalizable to all orders of perturbation theory . The Ginsparg-Wilson relation allows us to write, at finite lattice spacing, relations that are essentially equivalent to those holding in the low-energy phenomenology associated with chiral symmetry (see e.g. Refs. ). It indeed implies the existence of an exact chiral symmetry of the lattice action under the transformation
$$\delta _ϵ\psi (x)=ϵ\gamma _5(1aD)\psi (x),\delta _ϵ\overline{\psi }(x)=\overline{\psi }(x)\gamma _5ϵ,$$
(6)
which can be extended to the flavour non-singlet case<sup>*</sup><sup>*</sup>* Actually Ginsparg-Wilson fermions are invariant under a more general transformation
$$\delta _ϵ\psi (x)=ϵ\gamma _5(1vaD)\psi (x),\delta _ϵ\overline{\psi }(x)=\overline{\psi }(x)\left[1(1v)aD\right]\gamma _5ϵ,$$
(7) where $`v`$ can be chosen arbitrarily. All these symmetries are essentially equivalent, leading to equivalent Ward identities .. The axial anomaly then arises from the non-invariance of the fermion integral measure under flavour-singlet chiral transformations . A lattice formulation of QCD satisfying the Ginsparg-Wilson relation overcomes the complications of the standard approach (e.g. Wilson fermions), where chiral symmetry is violated at the scale of the lattice spacing. The chiral symmetry ensures that the hadron masses are free of $`O(a)`$ discretization errors. Indeed their leading scaling corrections are $`O(a^2)`$. The important point is that lattice Dirac operators satisfying Eq. (5) are not affected by the Nielsen-Ninomiya theorem , thus they need not suffer from fermion doubling.
Indeed, $`D_\mathrm{N}`$ avoids fermion doubling. The locality properties of $`D_\mathrm{N}`$ in the presence of a gauge field are not obvious. $`D_\mathrm{N}`$ is not strictly local, but locality should be recovered in a more general sense, i.e. allowing an exponential decay of the kernel of $`D_\mathrm{N}`$ with a rate which scales with the lattice spacing and not with the physical quantities . Thus, the Neuberger-Dirac operator seems to have all the right properties that a lattice Dirac operator should have in order to describe massless quarks. However, its complexity renders its numerical implementation very demanding. In this respect some progress has been achieved (see, e.g., Refs. ), and Monte Carlo simulations seem to be already feasible, at least in quenched approximation.
The matrix elements of the fermionic currents are employed to predict hadronic decay constants, electromagnetic and weak form factors, quark masses, etc. A knowledge of their lattice renormalization constants is necessary to relate the matrix elements computed using lattice simulations to the corresponding ones defined using continuous renormalization schemes in experimental data analysis. In this paper we compute, to one loop in perturbation theory, the renormalization constants of local bilinear quark operators $`\overline{\psi }\mathrm{\Gamma }\psi `$, where $`\mathrm{\Gamma }`$ denotes the generic Dirac matrix, in the lattice formulation of QCD using the Neuberger-Dirac operator. We also extend our computation to quark bilinears which are more extended and have improved chiral properties. In view of the better chiral behaviour of the Neuberger-Dirac operator, we expect it to have a wider use in Monte Carlo studies of hadronic physics; such studies require knowledge of the above renormalization constants.
Perturbative computations are much more cumbersome than in the case of Wilson fermions, due to the more complicated structure of the Neuberger-Dirac operator . One-loop calculations are already quite complicated, and require the use of symbolic manipulation packages. We have developed such a package in Mathematica. An example of one-loop calculation using the Neuberger-Dirac operator is reported in Ref. , where we computed the ratio $`\mathrm{\Lambda }_L/\mathrm{\Lambda }_{\overline{MS}}`$ between the $`\mathrm{\Lambda }`$-parameters of the lattice formulation and of the $`\overline{\mathrm{MS}}`$ renormalization scheme.
Nonperturbative methods to estimate the renormalization constants, such as those of Refs. would in general be preferable to approximations based on perturbative calculations, due to their better controlled systematic errors ($`O(a)`$ against $`O(g_0^n)`$). However, perturbative estimates may still be quite useful. They indeed provide important consistency checks. Further, in those cases where nonperturbative methods are very costly to implement, as in the case of the Neuberger formulation of the Dirac operator, perturbative methods may remain the only source of quantitative information.
One could, of course, define exactly conserved vector and axial currents having $`Z_A=Z_V=1`$, following Noether’s procedure, (i.e. performing variations of the action with respect to the non-singlet chiral transformations of Eq. (6)). However, these Noether currents are rather complicated extended objects, which might be cumbersome to use in actual simulations. Their expressions for the Neuberger-Dirac operator can be found in Ref. .
The paper is organized as follows: In Sec. II we report the one-loop calculations necessary for the evaluation of the lattice renormalizations of the two-quark operators. In Sec. III we improve the one-loop estimates by performing a resummation to all orders of a certain class of gauge invariant diagrams, dubbed “cactus” diagrams . We also compare them with the results of other improvement recipes , such as the so-called tadpole improved perturbation theory.
## II One-loop lattice renormalization of the two-quark operators
### A Formulation of the problem
The lattice regularization of QCD we consider is described by the action
$$S_\mathrm{L}=\frac{1}{g_0^2}\underset{x,\mu ,\nu }{}\mathrm{Tr}\left[1U_{\mu \nu }(x)\right]+\underset{i=1}{\overset{N_f}{}}\underset{x,y}{}\overline{\psi }_i(x)D_\mathrm{N}(x,y)\psi _i(y).$$
(8)
where $`U_{\mu \nu }(x)`$ is the usual product of $`SU(N)`$ link variables $`U_\mu (x)`$ along the perimeter of a plaquette originating at $`x`$ in the positive $`\mu `$-$`\nu `$ directions, and $`N_f`$ is the number of massless quark flavours.
The observables we study are bilinear quark operators of the form
$$O_i=\overline{\psi }(x)\mathrm{\Gamma }_i\psi (x),$$
(9)
where $`\mathrm{\Gamma }_i`$ denotes generic Dirac matrices, i.e.
$$1,\gamma _5,\gamma _\mu ,\gamma _\mu \gamma _5,\sigma _{\mu \nu }\gamma _5.$$
(10)
Specific bilinear operators are denoted according to their Lorentz group transformations: $`S(x)\overline{\psi }(x)\psi (x)`$, $`P(x)\overline{\psi }(x)\gamma _5\psi (x)`$, $`V_\mu (x)\overline{\psi }(x)\gamma _\mu \psi (x)`$, $`A_\mu (x)\overline{\psi }(x)\gamma _\mu \gamma _5\psi (x)`$, and $`T_{\mu \nu }(x)\overline{\psi }(x)\sigma _{\mu \nu }\gamma _5\psi (x)`$. In the above definitions, flavour indices are left unspecified. Indeed, this is sufficient for our purpose, since one-loop calculations are unaffected.
The lattice renormalization constants $`Z_{O_i}`$ of the operators $`O_i`$ are defined once a renormalization scheme is given: we will consider the $`\overline{\mathrm{MS}}`$ renormalization scheme. In order to simplify our calculations we work with zero quark mass. This is justified by the fact that renormalization constants in the $`\overline{\mathrm{MS}}`$ scheme are independent of the fermionic mass.
We also introduce the following bilinear operators
$`S^{}\overline{\psi }\left(1{\displaystyle \frac{1}{2}}aD_\mathrm{N}\right)\psi ,`$ $`P^{}\overline{\psi }\gamma _5\left(1{\displaystyle \frac{1}{2}}aD_\mathrm{N}\right)\psi ,`$ (11)
$`V_\mu ^{}\overline{\psi }\gamma _\mu \left(1{\displaystyle \frac{1}{2}}aD_\mathrm{N}\right)\psi ,`$ $`A_\mu ^{}\overline{\psi }\gamma _\mu \gamma _5\left(1{\displaystyle \frac{1}{2}}aD_\mathrm{N}\right)\psi ,`$ (12)
which are extended, and local only in the more general sense holding for $`D_\mathrm{N}`$. They are interesting since under the lattice chiral symmetry of Eq. (6) they transform as the corresponding operators in the continuum. Indeed one can easily check that
$$\delta _ϵS^{}=2P^{}ϵ,\delta _ϵP^{}=2S^{}ϵ,$$
(13)
and under the non-singlet transformations
$$\delta _{ϵ^a}\psi _i=ϵ^aT_{ij}^a\gamma _5(1aD)\psi _j,\delta _{ϵ^a}\overline{\psi }_i=\overline{\psi }_j\gamma _5T_{ji}^aϵ^a,$$
(14)
the vector and axial currents $`V^{}`$ and $`A^{}`$ transform as
$$\delta _{ϵ^b}V^a=if^{abc}ϵ^bA^c,\delta _{ϵ^b}A^a=if^{abc}ϵ^bV^c.$$
(15)
Using the corresponding Ward identities one can show that
$$Z_S^{}=Z_P^{},Z_V^{}=Z_A^{}.$$
(16)
We mention that, as pointed out in Refs. , the spontaneous breaking of the non-singlet chiral symmetry, Eq. (6), is related, for finite lattice spacing, to the vacuum expectation value (condensate) of the operator $`S^{}`$, i.e. it occurs if $`S^{}0`$.
One can also prove that the renormalization constants of $`O_i^{}\overline{\psi }\mathrm{\Gamma }_i(1\frac{1}{2}aD_N)\psi `$ coincide with those of the corresponding operators $`O_i\overline{\psi }\mathrm{\Gamma }_i\psi `$. This can be seen from the path integral representation for $`Z_{O_i}`$: Upon fermionic integration, the factor of $`D_N`$ in the definition of $`O_i^{}`$ will cancel against a fermion propagator; there will remain another propagator which will give a vanishing contribution to $`Z_{O_i}`$, since it cannot develop any $`𝒪(1/a)`$ singularities (by virtue of chiral invariance, which forbids additive mass renormalization). Thus, using the relations Eq. (16), we also have that
$$Z_S=Z_P,Z_V=Z_A.$$
(17)
Note that if one introduces the bare fermionic mass by writing the Dirac operator as in Ref.
$$\left(1\frac{1}{2}am_0\right)D_\mathrm{N}+m_0,$$
(18)
then the renormalization of the fermionic mass is related to that of the operator $`\overline{\psi }\left(1\frac{1}{2}aD_\mathrm{N}\right)\psi `$, denoted $`Z_S`$, by the relation
$$Z_m=Z_S^1.$$
(19)
It is worth mentioning that the chiral symmetry of the Neuberger-Dirac operator ensures the absence of $`O(a)`$ discretization errors only for the spectrum of the theory. In order to achieve the same property for generic matrix elements of local operators, one should employ improved operators . This issue is discussed in Ref. for bilinear operators constructed with Ginsparg-Wilson fermions; there it is shown that, for massless quarks, such an improvement can be achieved using the operators
$$O_i^{\prime \prime }\overline{\psi }\left(1\frac{1}{2}aD_\mathrm{N}\right)\mathrm{\Gamma }_i\left(1\frac{1}{2}aD_\mathrm{N}\right)\psi .$$
(20)
Proceeding as before, one can prove that the renormalization constants of the operators $`O_i^{\prime \prime }`$ are the same as those of the corresponding $`O_i`$.
The standard perturbative computation of the lattice renormalizations $`Z_{O_i}`$ of the operators $`O_i`$ requires the calculation of the two-quark one-particle irreducible function, $`\mathrm{\Gamma }(p)`$, and the two-quark one-particle irreducible functions $`\mathrm{\Gamma }_{O_i}(p)`$ with an insertion of the operators $`O_i`$. In the massless case, the quark-field renormalization $`Z_\psi `$ is obtained through the relation
$$\mathrm{\Gamma }^{\overline{\mathrm{MS}}}(p)=Z_\psi (a\mu )\mathrm{\Gamma }^\mathrm{L}(p)$$
(21)
where $`\mathrm{\Gamma }^\mathrm{L}(p)`$ is the two-quark function calculated on the lattice (in the limit $`a0`$), and $`\mathrm{\Gamma }^{\overline{\mathrm{MS}}}(p)`$ is the $`\overline{\mathrm{MS}}`$-renormalized two-point function, which can be computed in the continuum. Indeed a simple calculation gives
$$\mathrm{\Gamma }^{\overline{\mathrm{MS}}}(p)=ip_\mu \gamma _\mu \left[1+g^2c_F\left(c\mathrm{ln}\frac{\mu ^2}{p^2}+b^{\overline{\mathrm{MS}}}\right)+O(g^4)\right],$$
(22)
where
$$c=\frac{\alpha }{16\pi ^2},b^{\overline{\mathrm{MS}}}=\frac{\alpha }{16\pi ^2},$$
(23)
and $`\alpha `$ is the gauge parameter. The renormalizations $`Z_{O_i}`$ can be obtained by the equation
$$\mathrm{\Gamma }_{O_i}^{\overline{\mathrm{MS}}}(p)=Z_{O_i}(a\mu )Z_\psi (a\mu )\mathrm{\Gamma }_{O_i}^\mathrm{L}(p)$$
(24)
where $`\mathrm{\Gamma }_{O_i}^\mathrm{L}(p)`$ and $`\mathrm{\Gamma }_{O_i}^{\overline{\mathrm{MS}}}(p)`$ are the two-quark functions with an insertion of $`O_i`$ calculated respectively on the lattice and in the continuum using the $`\overline{\mathrm{MS}}`$ renormalization scheme. Setting
$$\mathrm{\Gamma }_{O_i}^{\overline{\mathrm{MS}},\mathrm{L}}(p)=\mathrm{\Gamma }_iB_{O_i}^{\overline{\mathrm{MS}},\mathrm{L}}(p),$$
(25)
one has
$$B_{O_i}^{\overline{\mathrm{MS}}}(p)=1+g^2c_F\left(c_{O_i}\mathrm{ln}\frac{\mu ^2}{p^2}+b_{O_i}^{\overline{\mathrm{MS}}}\right)+O(g^4),$$
(26)
where (setting $`\alpha =1`$)
$$c_{S,P}=\frac{1}{4\pi ^2},c_{V,A}=\frac{1}{16\pi ^2},c_T=0,$$
(27)
and
$$b_{S,P}^{\overline{\mathrm{MS}}}=\frac{3}{8\pi ^2},b_{V,A}^{\overline{\mathrm{MS}}}=\frac{1}{16\pi ^2},b_T^{\overline{\mathrm{MS}}}=0.$$
(28)
In the above formulae we have neglected terms proportional to $`p_\mu p_\nu /p^2`$; these are present also in the lattice expressions, and cancel out in the calculation of the renormalizations.
On the lattice one finds
$$\mathrm{\Gamma }^\mathrm{L}(p)=ip_\mu \gamma _\mu \left[1+g^2c_F\left(c\mathrm{ln}a^2p^2+b^\mathrm{L}\right)+O(g^4)\right],$$
(29)
and
$$B_{O_i}^\mathrm{L}(p)=1+g^2c_F\left(c_{O_i}\mathrm{ln}a^2p^2+b_{O_i}^\mathrm{L}\right)+O(g^4).$$
(30)
Using Eqs. (21) and (24), we obtain
$`Z_\psi `$ $`=`$ $`1+g^2c_F\left(c\mathrm{ln}a^2\mu ^2+b^{\overline{\mathrm{MS}}}b^\mathrm{L}\right)+O(g^4),`$ (31)
$`Z_{O_i}`$ $`=`$ $`1+g^2c_F\left[\left(c_{O_i}c\right)\mathrm{ln}a^2\mu ^2+b_{O_i}^{\overline{\mathrm{MS}}}b^{\overline{\mathrm{MS}}}b_{O_i}^\mathrm{L}+b^\mathrm{L}\right]+O(g^4),`$ (32)
Note that $`Z_{O_i}`$ are independent of the gauge parameter $`\alpha `$.
So, to compute $`Z_{O_i}`$ we need to evaluate the constants $`b^\mathrm{L}`$ and $`b_{O_i}^\mathrm{L}`$ by a one-loop perturbative calculation on the lattice. We have to expand the Neuberger-Dirac operator in powers of $`g_0`$, and use the resulting vertices to construct the diagrams related to the one-particle irreducible functions $`\mathrm{\Gamma }^\mathrm{L}(p)`$ and $`\mathrm{\Gamma }_{O_i}^\mathrm{L}(p)`$. In the appendix we list the relevant formulae for our one-loop calculations.
### B Results and discussion
Two diagrams contribute to $`\mathrm{\Gamma }^\mathrm{L}(p)`$, shown in Figure 1; for $`B_{O_i}^\mathrm{L}(p)`$, there is only one diagram to one-loop, shown in Figure 2. Given that the 4-point vertex contains a part with an internal momentum ($`k`$ in Eq. (A15)), the corresponding part of the second diagram in Fig. 1 actually has the same connectivity as the first diagram.
The algebra involving lattice quantities was performed using a symbolic manipulation package which we have developed in Mathematica. For the purposes of the present work, this package was augmented to include the propagator and vertices of the overlap action.
We express our results in the form of Eqs. (29), (30). The parameters $`b^\mathrm{L},b_{O_i}^\mathrm{L}`$ depend on $`\rho `$, but not on $`N`$ or $`N_f`$.
To extract the $`p`$-dependence, we first isolate the divergent terms; these are responsible for the logarithms. There are only a few such terms, and in the pure gluonic case their values are well known. We can use these values also in diagrams with fermions, applying successive subtractions of the type:
$$\frac{1}{\overline{q}^2}=\frac{1}{\widehat{q}^2}+\left(\frac{1}{\overline{q}^2}\frac{1}{\widehat{q}^2}\right)$$
(33)
where $`\overline{q}^2`$ is the inverse fermionic propagator. All remaining terms now contain no divergences, and can be evaluated by Taylor expansion in $`ap`$.
At this stage, one is left with expressions which no longer contain $`p`$ and must be numerically integrated over the loop momentum. Given the complicated form of the overlap vertices, these expressions turn out to be quite lengthy, containing a few hundred terms in the cases at hand.
The integration is done in momentum space over finite lattices; an extrapolation to infinite size is then performed, in the manner of Ref. . We evaluated the integrals for a range of values of the parameter $`\rho `$, as presented in Table 1. For all values of $`\rho `$ that we quote, lattice sizes $`L128`$ are sufficient to yield answers to at least 7 significant digits (the uncertainty coming from a systematic error in the extrapolation, which can be estimated quite accurately). As the endpoints of the perturbative domain of $`\rho `$ are approached ($`\rho 0,\rho 2`$), some of the quantities we calculate require increasingly larger lattices, for similar accuracy; this is, of course, a reflection of the divergences in the propagator at these endpoints. Figure 3 shows the dependence of our results on $`\rho `$, for the whole range $`0<\rho <2`$.
A number of consistency checks, some of which are rather nontrivial, may be performed on our results. In particular:
* The logarithmic coefficients must equal those of the continuum.
* Terms proportional to $`p_\mu p_\nu /p^2`$, which appear in $`Z_V,Z_A`$, should match those of the continuum.
* $`Z_V=Z_A,Z_S=Z_P`$.
* No $`𝒪(p^0)`$ terms must appear in $`\mathrm{\Gamma }^\mathrm{L}(p)`$, i. e. no additive mass renormalization, as required by chiral symmetry.
Our results fulfill all of the above requirements. The last one serves also to verify our estimates of the systematic errors coming from the extrapolation.
For quick reference, we write the one-loop values of $`Z_\psi (\alpha =1),Z_{O_i}`$ at $`\rho =1`$, as follows:
$`Z_\psi `$ $`=`$ $`1+g^2c_F\left[(\mathrm{ln}a^2\mu ^2)/16\pi ^20.231966\right]`$ (34)
$`Z_{S,P}`$ $`=`$ $`1+g^2c_F\left[\mathrm{\hspace{0.17em}3}(\mathrm{ln}a^2\mu ^2)/16\pi ^2+0.204977\right]`$ (35)
$`Z_{A,V}`$ $`=`$ $`1+g^2c_F\left[\mathrm{\hspace{0.17em}0.198206}\right]`$ (36)
$`Z_T`$ $`=`$ $`1+g^2c_F\left[(\mathrm{ln}a^2\mu ^2)/16\pi ^2+0.204392\right]`$ (37)
## III Improved perturbation theory
In order to improve the estimates coming from lattice perturbation theory, one may perform a resummation to all orders of the so-called “cactus” diagrams . Briefly stated, these are gauge–invariant tadpole diagrams which become disconnected if any one of their vertices is removed. The original motivation of this procedure is the well known observation of “tadpole dominance” in lattice perturbation theory. In the following we refer to Ref. for definitions and analytical results.
Since the contribution of standard tadpole diagrams is not gauge invariant, the class of gauge invariant diagrams we are considering needs further specification. By the Baker-Campbell-Hausdorff (BCH) formula, the product of link variables along the perimeter of a plaquette can be written as
$`U_{x,\mu \nu }`$ $`=e^{ig_0A_{x,\mu }}e^{ig_0A_{x+\mu ,\nu }}e^{ig_0A_{x+\nu ,\mu }}e^{ig_0A_{x,\nu }}`$ (40)
$`=\mathrm{exp}\left\{ig_0(A_{x,\mu }+A_{x+\mu ,\nu }A_{x+\nu ,\mu }A_{x,\nu })+𝒪(g_0^2)\right\}`$
$`=\mathrm{exp}\left\{ig_0F_{x,\mu \nu }^{(1)}+ig_0^2F_{x,\mu \nu }^{(2)}+𝒪(g_0^4)\right\}`$
The diagrams that we propose to resum to all orders are the cactus diagrams made of vertices containing $`F_{x,\mu \nu }^{(1)}`$. Terms of this type come from the pure gluon part of the lattice action. These diagrams dress the transverse gluon propagator $`P_A`$ leading to an improved propagator $`P_A^{(I)}`$, which is a multiple of the bare transverse one:
$$P_A^{(I)}=\frac{P_A}{1w(g_0)},$$
(41)
where the factor $`w(g_0)`$ will depend on $`g_0`$ and $`N`$, but not on the momentum. The function $`w(g_0)`$ can be extracted by an appropriate algebraic equation that has been derived in Ref. and that can be easily solved numerically; for $`SU(3)`$, $`w(g_0)`$ satisfies:
$$ue^{u/3}\left[u^2/34u+8\right]=2g_0^2,u(g_0)\frac{g_0^2}{4(1w(g_0))}.$$
(42)
The vertices coming from the gluon part of the action, Eq. (8), get also dressed using a procedure similar to the one leading to Eq. (41) . Vertices coming from the Neuberger-Dirac operator stay unchanged, since their definition contains no plaquettes on which to apply the linear BCH formula.
One can apply the resummation of cactus diagrams to the calculation of the renormalization of lattice operators. Approximate expressions are obtained by dressing the corresponding one-loop calculations. Applied to a number of cases of interest , this procedure yields remarkable improvements when compared with the available nonperturbative estimates. As regards numerical comparison with other improvement schemes, such as boosted perturbation theory , cactus resummation fares equally well on all the cases studied . It is worth mentioning in passing that cactus resummation also affords us a systematic means of improving perturbation theory, by successively dressing higher loop diagrams.
Let us consider the renormalizations $`Z_V,Z_A`$ of isovector fermionic currents, which are finite functions of the bare coupling $`g_0`$. At one-loop order we have
$$Z_{V,A}=1+g_0^2z_{V,A}+\mathrm{},$$
(43)
where the constants $`z_V`$ and $`z_A`$ have been calculated in Sec.II. The cactus dressing of the above one-loop expressions can be simply obtained by using the dressed transverse gluon propagator, Eq. (41). We thus obtain the following approximate expressions
$$Z_{V,A}1+g_0^2\frac{z_{V,A}}{1w(g_0^2)}$$
(44)
Let us apply the above formula to the lattice $`SU(3)`$ gauge theory, for which $`z_{V,A}=0.26427`$, at $`g_0=1`$ which is a typical value for Monte Carlo simulation. Since ,
$$1w(g_0=1)=0.749775,$$
(45)
we find $`Z_{V,A}(\rho =1)1.35`$
Other recipes of improvement have been proposed in the literature (see e.g. , and for a review of them) that essentially consist in a better choice of the expansion parameter. Among them we mention the so-called tadpole improvement (MFI) motivated by mean-field arguments, in which one scales the link variable with $`u_0(g_0^2)\frac{1}{N}\mathrm{Tr}U_{x,\mu \nu }^{1/4}`$ as measured in the Monte Carlo simulation. Accordingly one rescales the coupling constant: $`g_0^2g_{\mathrm{mf}}^2=g_0^2/u_0^4`$. Thus, one obtains a mean-field improved expansion
$$Z_{V,A}=u_0\left[1+g_{\mathrm{mf}}^2\left(z_{V,A}+\frac{1}{12}\right)+O(g_{\mathrm{mf}}^4)\right]$$
(46)
For example, for $`SU(3)`$ in the quenched approximation and at $`g_0^2=1`$ one finds $`u_00.878`$ and $`g_{\mathrm{mf}}^21.68`$. Putting these number in Eq. (46) we obtain $`Z_{V,A}1.39`$, which is in reasonable agreement with the estimate coming using the cactus resummation.
Acknowledgements: H. P. would like to acknowledge the warm hospitality extended to him by the Theory Group in Pisa during various stages of this work. H. P. and E.V. would like to thank L. Del Debbio for useful discussions.
## A
In order to perform the lattice perturbative calculation we must formally expand $`D_\mathrm{N}`$ in powers of $`g_0`$. We list here the relevant expressions for the propagator and vertices, following Ref. .
Let us first write down the weak coupling expansion of the Wilson-Dirac operator $`D_\mathrm{W}`$. This will be useful for constructing the relevant vertices of $`D_\mathrm{N}`$. We write
$$X(q,p)=D_\mathrm{W}(q,p)\frac{1}{a}\rho =X_0(p)(2\pi )^4\delta ^4(qp)+X_1(q,p)+X_2(q,p)+O(g_0^3),$$
(A1)
where
$$X_0(p)=\frac{i}{a}\underset{\mu }{}\gamma _\mu \mathrm{sin}ap_\mu +\frac{1}{a}\underset{\mu }{}(1\mathrm{cos}ap_\mu )\frac{1}{a}\rho ,$$
(A2)
$`X_1(q,p)`$ $`=`$ $`g_0{\displaystyle d^4k\delta (qpk)A_\mu (k)V_{1,\mu }(p+k/2)},`$ (A3)
$`V_{1,\mu }(q)`$ $`=`$ $`i\gamma _\mu \mathrm{cos}aq_\mu +\mathrm{sin}aq_\mu ,`$ (A4)
$`X_2(q,p)`$ $`=`$ $`{\displaystyle \frac{g_0^2}{2}}{\displaystyle \frac{d^4k_1d^4k_2}{(2\pi )^4}\delta (qpk_1k_2)A_\mu (k_1)A_\mu (k_2)V_{2,\mu }(p+k_1/2+k_2/2)},`$ (A5)
$`V_{2,\mu }(q)`$ $`=`$ $`i\gamma _\mu a\mathrm{sin}aq_\mu +a\mathrm{cos}aq_\mu .`$ (A6)
The Fourier transform of the Neuberger-Dirac operator takes the form
$$\frac{1}{\rho }D_\mathrm{N}(q,p)=D_0(p)(2\pi )^4\delta ^4(qp)+\mathrm{\Sigma }(q,p).$$
(A7)
$`D_0(p)`$ is the tree level inverse propagator:
$$D_0^1(p)=\frac{i\underset{\mu }{}\gamma _\mu \mathrm{sin}ap_\mu }{2\left[\omega (p)+b(p)\right]}+\frac{a}{2},$$
(A8)
where
$`\omega (p)`$ $`=`$ $`{\displaystyle \frac{1}{a}}\left({\displaystyle \underset{\mu }{}}\mathrm{sin}^2ap_\mu +\left[{\displaystyle \underset{\mu }{}}(1\mathrm{cos}ap_\mu )\rho \right]^2\right)^{1/2},`$ (A9)
$`b(p)`$ $`=`$ $`{\displaystyle \frac{1}{a}}{\displaystyle \underset{\mu }{}}(1\mathrm{cos}ap_\mu ){\displaystyle \frac{1}{a}}\rho .`$ (A10)
The function $`\mathrm{\Sigma }(q,p)`$ can be expanded in powers of $`g_0`$ as
$`a\mathrm{\Sigma }(q,p)=`$ $`{\displaystyle \frac{1}{\omega (p)+\omega (q)}}\left[X_1(q,p){\displaystyle \frac{1}{\omega (p)\omega (q)}}X_0(p)X_1^{}(p,q)X_0(q)\right]`$ (A15)
$`+{\displaystyle \frac{1}{\omega (p)+\omega (q)}}\left[X_2(q,p){\displaystyle \frac{1}{\omega (p)\omega (q)}}X_0(p)X_2^{}(p,q)X_0(q)\right]`$
$`+{\displaystyle }{\displaystyle \frac{d^4k}{(2\pi )^4}}{\displaystyle \frac{1}{\omega (p)+\omega (q)}}{\displaystyle \frac{1}{\omega (p)+\omega (k)}}{\displaystyle \frac{1}{\omega (q)+\omega (k)}}\times `$
$`[X_0(p)X_1^{}(p,k)X_1(k,q)X_1(p,k)X_0^{}(k)X_1(k,q)X_1(p,k)X_1^{}(k,q)X_0(q)`$
$`+{\displaystyle \frac{\omega (p)+\omega (q)+\omega (k)}{\omega (p)\omega (q)\omega (k)}}X_0(p)X_1^{}(p,k)X_0(k)X_1^{}(k,q)X_0(q)]+\mathrm{}`$
From $`\mathrm{\Sigma }(q,p)`$ one can read off the vertices necessary for the one-loop calculations presented in this paper. |
warning/0002/hep-ph0002145.html | ar5iv | text | # Some remarks on the interpretation of degree of nonextensivity
## Abstract
Recently we have demostrated that the nonextensitivity parameter $`q`$ occuring in some applications of Tsallis statistics (known also as index of the corresponding Lévy distribution) is, in $`q>1`$ case, given entirely by the fluctuations of the parameters of the usual exponential distribution. We show here that this interpretation is valid also for the $`q<1`$ case. The parameter $`q`$ is therefore a measure of fluctuations of the parameters of the usual exponential distribution.
PACS numbers: 05.40.Fb 24.60.-k 05.10.Gg
Keywords: Nonextensive statistics, Lévy distributions, Thermal models
Recently we have demostrated that the nonextensitivity parameter $`q`$ occuring in some applications of Tsallis statistics (known also as index of the corresponding Lévy distribution) is, in $`q>1`$ case, given entirely by the fluctuations of the parameters of the usual exponential distribution . It means that:
* when in some exponential formula describing distribution of a quantity $`\epsilon `$ of physical interest:
$$L_{q=1}(\epsilon )=C_{q=1}\mathrm{exp}\left[\frac{\epsilon }{\chi }\right],$$
(1)
one allows the parameter $`\chi `$ to fluctuate around some mean value $`\chi _0`$, and
* if these fluctuations are described by simple Gamma distribution of the form
$$f(\frac{1}{\chi })=\frac{1}{\mathrm{\Gamma }(\alpha )}\mu \left(\frac{\mu }{\chi }\right)^{\alpha 1}\mathrm{exp}\left(\frac{\mu }{\chi }\right)$$
(2)
depending on two parameters
$$\alpha =\frac{1}{q1}\mathrm{and}\mu =\alpha \chi _0,$$
(3)
* then, as result, one gets the following power-like distribution for the quantity $`\epsilon `$ of interest:
$$L_q(\epsilon )=C_q\left[1(1q)\frac{\epsilon }{\chi _0}\right]^{\frac{1}{1q}},$$
(4)
known also as Lévy distribution with index $`q`$, where
$$q=\mathrm{\hspace{0.17em}1}+\frac{\left(\frac{1}{\chi }\right)^2\frac{1}{\chi }^2}{\frac{1}{\chi }^2},$$
(5)
i.e., where it is entirely given by the relative variance of the parameter $`1/\chi `$ of the initial distribution (1) ($`<\mathrm{}>`$ denotes the corresponding averages with respect to distribution $`f(\chi )`$).
The proof presented in was limited to the $`q>1`$ case and the physical discussion provided there was also concentrated on such situation. Because $`q`$ can be also interpreted as the so called nonextensivity parameter occuring in some applications of Tsallis statistic , it would be interesting to check if such interpretation can be extended to the $`q<1`$ case as well. We shall demonstrate below that this is indeed the case .
The essential difference between these two cases is, for the purpose of present discussion, that whereas for $`q>1`$ probability distribution $`L_q(\epsilon )`$ is well defined for the whole range of variable $`\epsilon `$, $`\epsilon (0,\mathrm{})`$, for $`q<1`$ it is defined only for $`\epsilon [0,\chi _0/(1q)]`$. As it was done in we shall deduce the form of function $`f(1/\chi )`$, describing fluctuations in $`\chi `$, which would lead from the exponential distribution $`L_{q=1}`$ to the power-like Lévy distribution $`L_{q<1}`$
$$L_{q<1}(\epsilon ;\chi _0)=C_q\left[1\frac{\epsilon }{\alpha ^{}\chi _0}\right]^\alpha ^{}=C_q_0^{\mathrm{}}\mathrm{exp}\left(\frac{\epsilon }{\chi }\right)f\left(\frac{1}{\chi }\right)d\left(\frac{1}{\chi }\right)$$
(6)
(for simplicity we denote $`\alpha ^{}=\frac{1}{1q}`$). From the representation of the Euler gamma function we have
$$\left[1\frac{\epsilon }{\alpha ^{}\chi _0}\right]^\alpha ^{}=\left(\frac{\alpha ^{}\chi _0}{\alpha ^{}\chi _0\epsilon }\right)^\alpha ^{}=\frac{1}{\mathrm{\Gamma }(\alpha ^{})}_0^{\mathrm{}}𝑑\eta \eta ^{\alpha ^{}1}\mathrm{exp}\left[\eta \left(1+\frac{\epsilon }{\alpha ^{}\chi _0\epsilon }\right)\right].$$
(7)
Changing now variables under the integral in such a way that $`\chi =\frac{\alpha ^{}\chi _0\epsilon }{\eta }`$ one immediately obtains Eq. (6) with $`f(1/\chi )`$ given by the following gamma distribution
$$f\left(\frac{1}{\chi }\right)=\frac{1}{\mathrm{\Gamma }(\alpha ^{})}\left(\alpha ^{}\chi _0\epsilon \right)\left(\frac{\alpha ^{}\chi _0\epsilon }{\chi }\right)^{\alpha ^{}1}\mathrm{exp}\left(\frac{\alpha ^{}\chi _0\epsilon }{\chi }\right)$$
(8)
with parameters $`\alpha ^{}`$ and $`\mu (\epsilon )=\alpha ^{}\chi _0\epsilon `$. This time, contrary to the $`q>1`$ case of , fluctuations depend on the value of the variable in question, i.e., the mean value and variance are both $`\epsilon `$-dependent:
$$\frac{1}{\chi }=\frac{1}{\chi _0\frac{\epsilon }{\alpha ^{}}}\mathrm{and}\left(\frac{1}{\chi }\right)^2\frac{1}{\chi }^2=\frac{1}{\alpha ^{}}\frac{1}{\left(\chi _0\frac{\epsilon }{\alpha ^{}}\right)^2}.$$
(9)
However, the relative variance
$$\omega =\frac{\left(\frac{1}{\chi }\right)^2\frac{1}{\chi }^2}{\frac{1}{\chi }^2}=\frac{1}{\alpha ^{}}=\mathrm{\hspace{0.17em}1}q$$
(10)
remains $`\epsilon `$-independent (exactly like in the case of $`q>1`$) and depends only on parameter $`q`$. It means therefore that the parameter $`q`$ in Lévy distribution $`L_q(\epsilon )`$ describes the relative variance of fluctuations of parameter $`\chi `$ in $`L_{q=1}(\epsilon )`$ for all values of $`q`$ (both for $`q>1`$, where $`\omega =q1`$, cf. and for $`q<1`$ as given above, where $`\omega =1q`$).
In we have proposed a general explanation of the meaning of function $`f(\chi )`$ describing fluctuations of some variable $`\chi `$. The question one is interested in is why, and under what circumstances, it is the gamma distribution that describes fluctuations. To this end we have started with general Langevin type equation for the variable $`\chi `$
$$\frac{d\chi }{dt}+\left[\frac{1}{\tau }+\xi (t)\right]\chi =\varphi =\mathrm{const}>\mathrm{\hspace{0.17em}0}$$
(11)
(with damping constant $`\tau `$ and source term $`\varphi `$). For stochastic processes defined by the white gaussian noise form of $`\xi (t)`$ (cf. for details) it can be shown that distribution function for the variable $`\chi `$ satisfies the Fokker-Planck equation ($`K_{1,2}`$ are the corresponding intensity coefficients, cf. )
$$\frac{df(\chi )}{dt}=\frac{}{\chi }K_1f(\chi )+\frac{1}{2}\frac{^2}{\chi ^2}K_2f(\chi ),$$
(12)
i.e., it is indeed given by the Gamma distribution in variable $`1/\chi `$ of the form (2) with $`\mu =\alpha \chi _0`$. Notice that it differs from Eq. (8) only in the form of parameter $`\mu `$, which in (8) depends also on the physical quantity of interest $`\epsilon `$.
As an illustration of the genesis of Eq. (11) we have discussed in the case of fluctuations of temperature (i.e., the situation where $`\chi =T`$) . Suppose that we have a thermodynamic system, in a small (mentally separated) part of which the temperature fluctuates around some mean value $`T_0`$ (which can be also understood as an equilibrium temperature) with $`\mathrm{\Delta }TT`$. The unevitable exchange of heat between this selected region and the rest of the system is described by Eq. (11) in which
$$\varphi =\varphi _{q<1}=\frac{1}{\tau }\left(T_0\frac{\epsilon }{\alpha ^{}}\right)\mathrm{whereas}\varphi =\varphi _{q>1}=\frac{T_0}{\tau }.$$
(13)
It means that the corresponding process of heat conductivity is, for $`q<1`$ case, described by the following equation (here $`T^{}=T_0\tau \xi (t)T`$)
$$\frac{T}{t}\frac{1}{\tau }(T^{}T)+\frac{\epsilon }{\tau \alpha ^{}}=\mathrm{\hspace{0.17em}0},$$
(14)
which differs from the corresponding equation for $`q>1`$ case only by the last term describing the presence the internal heat source. It has a sense of dissipative transfer of energy from the region where (due to fluctuation) we have higher $`T`$. It could be any kind of convection type flow of energy, for example it could be connected with emission of particles from that region. The heat release given by $`\epsilon /(\tau \alpha ^{})`$ depends on $`\epsilon `$ (but it is only a part of $`\epsilon `$, which is released). In the case of such energy release (connected with emission of particles) there is additionale cooling of the whole system. If this process is sufficiently fast, it could happen that there is no way to reach a stationary distribution of temperature (because the transfer of the heat from the outside can be not sufficient for development of the state of equilibrium). On the other hand (albeit this is not our case here) for the reverse process we could face the ”heat explosion” situation (which could happen if the velocity of the exotermic burning reaction grows sufficiently fast; in this case because of nonexistence of stationary distribution we have fast nonstationary heating of the substance and acceleration of the respective reaction).
It should ne noticed that in the case of $`q<1`$ the temperature does not reach stationary state because, cf. Eq. (9), $`1/T=\mathrm{\hspace{0.17em}1}/(T_0\epsilon /\alpha ^{})`$, whereas for $`q>1`$ we had $`<T>=T_0`$. As a consequence the corresponding Lévy distribution are defined only for $`\epsilon (0,T_0\alpha ^{}`$) because for $`\epsilon T_0\alpha ^{}`$ the $`<T>0`$. Such asymptotic (i.e., for $`t/\tau \mathrm{}`$) cooling of the system ($`T0`$) can be also deduced form Eq. (14) for $`\epsilon T_0\alpha ^{}`$.
To summarize, we have demonstrated that temperature fluctuations lead to the Lévy distribution $`L_q(\epsilon )`$ with index $`q<1`$ when there exists energy source and with $`q>1`$ in the absence of such source. In both cases, however, the relative variance of $`1/T`$ fluctuations is described by the parameter $`q`$ only. |
warning/0002/physics0002019.html | ar5iv | text | # Untitled Document
ON M THEORY, QUANTUM PARADOXES AND THE NEW RELATIVITY
Carlos Castro
Center for Theoretical Studies of Physical Systems
Clark Atlanta University
Atlanta, GA. 30314
Alex Granik
Department of Physics
University of the Pacific
Stockton, California 95211
January 2000
(Dedicated to the memory of Leonard Ainsworth, a true pioneer and a scientist )
Abstract
Recently a New Relativity Principle has been proposed by one of the authors as the underlying physical and geometrical foundations of String and M Theory. It is explicitly shown that within the framework of the New Relativity Theory, some Quantum Mechanical Paradoxes like the Einstein-Rosen Podolsky and the Black Hole Information Loss, are easily resolved. Such New Relativity Theory requires the introduction of an Infinite Dimensional Quantum Spacetime as has been shown recently by one of us. This can be viewed as just another way of looking at Feynman’s path integral formulation of Quantum Mechanics. Instead of having an infinite dimensional funcional integral over $`all`$ paths, smooth, forwards and backwards in time, random and fractal, in a finite-dimensional spacetime, one has a finite number of paths in an Infinite Dimensional Quantum Spacetime. We present a few-lines proof why there is no such a thing as an EPR Paradox in this New Relativity theory. The reason is not due to a superluminal information speed but to a divergent information charge density. In the infinite dimensional limit, due to the properties of gamma functions, the hypervolume enclosed by a $`D`$-dim hypersphere, of finite nonzero radius, shrinks to zero : to a hyperpoint , the infinite-dimensional analog af a point. For this reason, Information flows through the infinite-dimensional hypersurface of nonzero radius, but zero size, the hyperpoint, in an instant. In this fashion we imbue an abstract mathematical ”point” with a true physical meaning : it is an entity in infinite dimensions that has zero hypervolume at nonzero radius . A plausible resolution of the Information Loss Paradox in Black Holes is proposed.
1. Introduction : Historical background
From the very beginning the relativity principle has been one of the cornerstones of mechanics. In 1686 , in the opening pages of $`Principia`$, Newton wrote : ” I do not define time, space, place, and motion, as being well known to all. Only I must observe, that the common people conceive those quantities under no other notions but from the relations they bear to sensible objects ”. Based on this preamble, Newton introduced the absolute space, time, and motion which in his own words do not bear ” relation to anything external ” . It is clear that such absolute space, time, motion, are purely metaphysical notions, stand outside the realm of physics, and serve the purpose of ” geometrization ” of physical phenomena. Thus at the foundations of Newtonian mechanics lie the above metaphyisical concepts. The very same Metaphysical concepts that so many members of the scientific community, by their own admission, are unwilling to accept.
Much later, Mach rather broadly , and then Einstein definitely, set Physics on physical ground by defining measurements relative only to the physical phenomena, and not to the metaphysical entities. Instead of vaguely-defined ( if at all) metaphyiscal concepts of absolute space, time, and motion, new ( rather narrow) but well defined physical concepts of physical measurements had been introduced. The old ”absolutes ” had been dismissed by Einstein’s Theory of Relativity. Measurements required a universal standard which had been furnished by a physical quantity : the speed of light.
Based on this historical perspective we introduce the New Relativity Principle that rest on the four postulates presented in the next section. In the last sections an amazingly simple proof is presented why there is no EPR Paradox in such New Relativity Theory. In addition, a plausible resolution of the Information Loss in Black Holes is proposed.
2. The Program of the New Relativity Principle : The Demolition of Today’s Absolutes
Recently one of the authors has proposed that a New Relativity principle may be operating in Nature which could reveal important clues to find the origins of $`M`$ theory . We were forced to introduce this new Relativity principle, where all dimensions and signatures of spacetime are on the same footing, to find a fully covariant formulation of the $`p`$-brane Quantum Mechanical Loop Wave equations. This New Relativity Principle, or the principle of Polydimensional Covariance as has been called by Pezzaglia, has also been crucial in the derivation of Papapetrou’s equations of motion of a spinning particle in curved spaces that was a long standing problem which lasted almost 50 years . A Clifford calculus was used where all the equations were written in terms of Clifford-valued multivector quantities; i.e one had to abandon the use of vectors and tensors and replace them by Clifford-algebra valued quantities, matrices, for example .
The New Relativity Theory rests on four postulates :
1. The old Bootstrap Idea of Chew : Each $`p`$-brane is made of all the others. To view a single $`p`$-brane as a fundamental identity is a meaningless concept. $`p`$-branes are defined only in relation to others. This is Mach’s principle once again. For this reason, one must include all dimensions and signatures on the same footing. Pezzaglia has called this the principle of Polydimensional Covariance or Dimensional Democracy. The New Relativity theory reshuffles, for example, a string history for a $`5`$-brane history; a $`9`$-brane history for a membrane history; an $`11`$-brane history for a history encompassing all other $`p`$-branes…and so forth. Point and extended instantons and tensionless $`p`$-branes are also inlcluded. The tensionless $`p`$-brane history excitations of the infinite dimensional Quantum Spacetime are the ” photons ” in this New Relativity Theory.
We honestly believe that $`M`$ Theory does not stand for mystery, membrane, matrix, master, mother, murky, Moyal….it stands for Mach. The New Relativity theory is based on the ultimate Machian view of the Quantum Universe ( the Ultimate Machian ” Quantum Computer ” ) : Relationships among entities are the only meaningful statements one can make . A perfect example of this are : spin-networks, quantum-networks, quantum sets, cellular networks, $`p`$-Adic Physics,…etc. Since it is undesirable to run off the letters of the alphabet, by keep adding letters like $`M,F,S\mathrm{}`$ Theory, we gather courage to say that by abandoning the Egocentric Anthropomorphic view of the Universe and, instead, embracing Mach’s view that everything in the Quantum Universe is interconnected, one reaches the end of the alphabet at $`Z`$ theory : $`Z`$ stands for the ultimate Machian view of the Quantum Universe ( that dismisses the egocentric view of the Universe ) for a $`Zenthropic`$ Quantum Universe. Who are we to say that we know everything that an electron , a photon, a quark really ” sees ” ? Do electrons, quarks….perform Feynman diagrams ? Has anyone seen a point ?
2. Laurent Nottale’s Scale-Relativity theory .
In the final analysis, Physics involves, and is about, measurements. Physics is an experimental science. Physics deals with experiences. On the other hand to measure something one needs a standard of measurement to compare measurements with. It is essential , it is of prime importance, to introduce resolutions in Physics. It is meaningless to say that the one has a field at x. By x meaning : specifying the value of the real number $`x`$ to an infinite number of non-periodic decimal places. In mathematics we can infinitely increase the accuracy ( or degree of resolution) at will of any real number by adding digits. However, in practice we cannot have such arbitrary accuracy provided by mathematical constructions. Even writing an infinite aperiodic decimal fraction would require an infinite amount of memory. Therefore, in Physics, it is necessary to have a finite universal physical ”yardstick” which would define the ultimate Physical Resolution. Nottale’s Scale Realtivity takes the Planck scale as such Universal physical standard of measurement that is invariant, by definition, under Scale Relativistic transformations of resolutions, like the speed of light was in Einstein’s Relativity.
$`p`$-adic numbers and $`p`$-adic Physics is a nice attempt to eliminate the problem of having to specify a real number up to infinite digits . The Planck scale is therefore taken as that universal standard of measure invariant under Scale Relativistic transformations of resolutions . In the same vein that the speed of light was taken by Einstein as the maximum speed in Nature, the Planck scale is taken to be the minimum length. The speed of light allowed Einstein to embrace space with time, since space and time have different units. By the same token, to embrace all dimensions one needs a Universal length scale in all dimensions : the Planck Scale.
As the years pass by, more and more planets have been found confirming Nottale’s predictions within his framework of Scale-Relativity. Instead of being properly rewarded with increased curiosity and interest in his remarkable theory, he has been increasingly rewarded with insults and a suffocating censorship . As the number of his planet confirmations increases, so does the number of insults increases and the censorship of his work is tightened further. Unfortunately, the New Relativity Theory will never be able to explain such odd phenomena.
The Universal scale, in units of $`\mathrm{}=c=1`$, in any dimensions, $`D>2`$ ( in two dimensions the Einstein-Hilbert action is a topological invariant) is :
$$\mathrm{\Lambda }=G_D^{\frac{1}{D2}}=G_{D1}^{\frac{1}{D3}}=\mathrm{}\mathrm{}\mathrm{}..=10^{33}cms.$$
$`(1a)`$
where $`G_D,G_{D1}\mathrm{}..`$ are the Newton gravitational coupling constants in different dimensions. In the same fashion that in Newtonian Physics one only can assign a definite meaning to the ratio of masses ( it is meaningless to say that one has a value of $`m`$ without a comparison to another mass ), in the New Relativity theory it only makes sense to write :
$$\frac{D3}{D2}=\frac{lnG_{D1}}{lnG_D}.$$
$`(1b)`$
In the New Relativity Theory it is meaningless to talk about such things as ” compactification ” , ” decompactification ” used in the literature that relates the Newton constants in different dimensions through the small radius of a compactified unseen dimension at low energies. It as meaningless as saying that the velocities of the gas molecules in a room experience a dynamical or spontaneous ” compactification ” to a fixed average value.
Problems with the compactification picture of Superstring theory from 10 to 4 were already alarming signals that something could be wrong. Billions and billions of possible four-dimensional phenomenological theories of the world were obtained : the so-called uniqueness of string theory went out the window when this was found. String theory wasn’t the problem, assuming a fixed dimensions was ! Witten already proved long ago that something might be inherently wrong with the compactification schemes, when he showed using Index Theory arguments, that the standard 11-dim Supegravity Kaluza-Klein compactifications of ordinary manifolds did not yield chiral fermions in 4 dimensions. This problem was bypassed in the second string revolution by saying that orbifold compactification were fine because orbifolds are not really ordinary manifolds, so things were satisfactory after all. The New Relativity Theory does not have to face these challenges. One has a truly infinite-dimensional Quantum Universe which suggests that Topological Field Theories could be the most natural candidates for a theory of the world. Since below the Planck scale there is no such thing as a distance; it is very likely that Topology should play a more important role.
Conformal Field Theories and their Higher Conformal spin extensions are the ones to use in $`D=2`$. In $`D=2`$, one has induced gravity : $`W_2,W_3\mathrm{}W_{\mathrm{}}`$ gravity as a result of integrating out the conformal matter field fluctuations. This replaces the topological invariant Einstein-Hilbert action. In $`D=1`$ dimension there is only extrinsic curvature. One can view a one-dim loop as the boundary of a two-dimensional surface. This allowed to write down a String Representation of Quantum Loops from a covariantized phase space Schild action path integral. The effective action for the boundary, with induced extrinsic curvature terms was obtained, in addition to the Polyakov Bulk partition function and the holographic boundary Eguchi wave functional as well.
3. Noncommutative C-spaces. One of the authors was forced to enlarge the naive notion of commuting spacetime coordinates to fully covariantize the Quantum Mechanical Loop Equations for $`p`$-branes. One achieved that goal if one extended the notion of ordinary spacetime vectors and tensors, to a Noncommutative C-space, or Clifford manifold, where all $`p`$-branes were unified in one single footing by using Clifford-algebra valued multivectors quantities ( matrices) instead of ordinary vectors and tensors. In order to combine objects of different dimensionality one needs a length scale : the Planck scale.
There was a one-to-one correspondence between a nested hierarchy of point, loop, 2-loop, 3-loop …….p-loop histories in $`D`$ dimensions encoded in terms of hyper-matrices and single lines in Clifford Manifolds. This is roughly similiar to the aim of Penrose’s twistor progam. By using Clifford-algebra valued multivectors, one could argue why it may be meaningless to say that the cosmological constant is a constant in its definition ! The so-called cosmological constant is observer-dependent in this New Relativity Theory : it is just one of the many components of the Clifford multivectors. Due to Polydimensional Covariance, only the norm of such multivector is truly an invariant. So using this simple argument one of us was able to argue why it is meaningless to try to measure such constant, unless one is specifying what is the frame of reference one lives in !
The reader may say that the value of $`p=1`$ was not included here. Point and Extended Instantons can also be treated very naturally in this framework . The New Relativity Theory reshuffles, for example, a loop-history represented by the coordinates : $`x^\mu ,\sigma _{\mu \nu },A`$ in one frame of reference, to another history, in another frame of reference, represented by the loop-instanton $`x^\mu ,(\sigma _{\mu \nu })^{},A^{}=0`$. The $`x_{CM}`$ are the center of mass coordinates of the loop. $`A`$ is the areal-time spanned by the motion of the loop through spacetime. $`\sigma _{\mu \nu }`$ are the holographic coordinates of the loop. It can reshuffle a massive point history ( a line ) : $`x^\mu ,\tau 0`$ to a massive point-instanton : $`x^\mu ,\tau ^{}=0`$ in another frame of reference. An so forth.
4. Quantum Spacetime must be treated from a Multivector-Multiscale point of view. The use of Clifford-valued multivectors was explained above. The multi-scale or resolution aspects are based on Nottale’s fractals and El Naschie’ s Cantorian-Fractal Quantum Spacetime views that dimensions are resolution dependent concepts and not fixed notions .
Nottale, by abandoning the hypothesis of the differentiablity of spacetime , was led to three effects ( at least ) : (i) . The number of geodesics becomes infinite. This forces upon us to jump to a statistical fluid-like description. (ii) Each geodesic becomes a fractal curve of higher and higher fractal dimensionality as the resolution of the ” physical apparatus ” becomes finer and finer, asympotically approaching the minimum Planck scale resolution where the fractal dimensionality becomes infinite. This forces us to embed the fractal geodesics in an spacetime of infinite-Hausdorff dimensions. (iii). The symmetry $`dtdt`$ is broken by the non-differentiablity which leads to a two-valuedness character of the average velocity vector and which is, in Nottale’s view, the underlying reason why the wave function in QM is complex.
This is not the ultimate status of things. To be consistent and to move forward along the path charted by Mach and Einstein, one cannot, and should not , accept this status quo as the ” end of the road ” in Physics. This reminds us of the status of things at the end of the 19 century when ” two clouds ” were the only obstacles hovering over the horizon that prevented the ” end ” of Physics. In fact, one cannot but to feel compelled to say that from the beginning, a truly quantum mechanical description of the world must start by abandoning the $`verynotionofspacetimeitself`$ and other ” idols ” from our minds, as Finkesltein has pointed out . This is precisely the goal of $`p`$-Adic Physics to remove the notion of spacetime per se and replace it by objects and their relationships. A truly Categorical view of the Universe. An extension of Einstein’s motion Relativity and Nottale’s Scale Relativity into a unfied Scale-Motion Realtivity was outlined briefly in . Whatever the ”final” view of the world may be, it seems that it is wrong to assume that Quantum Spacetime has a fixed dimension. On the contrary, it may have uncountably-infinite dimensions as El Naschie has argued . Taking this infinite-dimensional point of view allows us to eliminate the notion of a EPR , and possibly, Black-Hole Information Loss ”paradoxes ”. For this reason we believe it ought to be investigated further. Dimensions are not fixed absolutes. They are resolution dependent concepts.
Quantum Gravity is not a quantization of the spacetime coordinates, metric…..If this were the case, one would have had quantized the spacetime coordinates long ago. In String Theory, from the two-dim world sheet point of view , the spacetime coordinates are nothing but a finite number of scalar fields whose quantization is essentially trivial by selecting the conformal or orthonormal gauge. The same arguments applies with the ( linearized ) spin two graviton. Quantum Gravity it is something much deeper than the naive notion of coordinates and gravitons. It is something that doesn’t need any spacetime background nor metrics whatsoever. Morever, it involves something that disposes of the ill-conceived notion of having a fixed dimension. The classical spacetime that we perceive with our senses is just a long distance averaging effect associated with a quantum network of processeses of a deeper underlying Quantum Universe. Einstein’s Gravity is an effective theory as suspected long ago. To merge Quantum Mechanics with Relativity it is necessary to enlarge the Einsteinian view of Relativity to a New Relativity Principle . To proceed further one has to demolish the concept of dimension as an absolute , as an idol.
To sum up what has been said so far :
The New Relativity Theory forces upon us to take a radically different view of the Quantum World, an ultimate Machian/Zenthropic view, and to dismiss the concepts of false absolutes (idols) of dimensions, spacetime, cosmological constant , from our classical minds, as Finkelstein has advocated. If a true evolution ( revolution) of Physics is to take place one has to embrace the plausible extensions of Relativity as Finkelstein has insisted . For those who believe that we have reached the end of the road, the end of Physics, we feel that they are setting themselves for similar surprises that Lord Kelvin experienced with the advent of Quantum Mechanics and Relativity.
To this day , to the best of our knowledge, there is no satisfactory definition yet of Quantum Field Theory. QFT today is being challenged by deeper concepts : Noncommutative Geometry, Quantum Groups, Hopf algebras, Monoidal Braided Categories, Braided QFT, etc…. Relativity itself is hereby extended to a deeper meaning by the New Relativity Theory : Scale Relativity and Cantorian-Fractal Geometry. A Nested Hierarchy of Histories have replaced the old fashioned concepts of spacetime events; vectors and tensors have been replaced by Clifford-multivectors; Riemannian Geometry by Finsler Geometry and by Fractal-Cantorian, Non-Archimedean, $`pAdic`$ , Noncommutative and Nonassociative Geometries…..
Recently we have proposed to even abandon the the idea of the cosmological constant as a constant. The so-called cosmological constant is not a constant in its definition ! It is observer-dependent within the framework of the New Relativity Theory . Trying to estimate the absolute values of such a ” constant ” is like trying to detect absolute spacetime motion and to verify the existence of the ether ! Such ideas that the vacuum energy could be observer dependent orginated with discussions held in Trieste by one of us with Miguel Cardenas and Devashis Benarjee . Even the notion of the ”vacuum” per se ! Special Relativity demolished such framework of thinking. We believe that the New Relativity Theory will also replace the existence of such ill-conceived notions that spacetime has a fixed dimension and that the cosmological constant has an well defined absolute value in all frame of references.
The observed spacetime dimension of $`D=4`$ is interpreted in this New Relativity Theory as a result of an averaging procedure over all the possible infinite values of Quantum Spacetime. In a sense it is similar to what happens with the statistical distribution of velocities of a gas. There is an average velocity ( average over all the infinite possible values of the statistical ensemble ) proportional to the Temperature. To assume that there is a spacetime compactification from D=11 to D=4 ( like it is assumed in mainstream Physics today ) is an incongrous assumption in this New Relativity Theory : it is like saying that there is a ” velocity compactification/decompactification ” from higher/lower velocities to the average observed velocity in a gas. Problems with the compactification picture of Superstring theory from 10 to 4 were already alarming signals when billions and billions of possible four-dimensional theories of the world were obtained : the so-called uniqueness of string theory went out the window when this was found. String theory wasn’t the problem, assuming a fixed dimensions was ! The fact that there might be an Statistical approach to the Dimensions, and to Quantum Gravity per se, was already lurking behind the scenes long ago in the work of Hawking : Black Hole Thermodynamics !
3. There is No Such Thing as an Einstein-Rosen-Podolski Paradox in the New Relativity
We will present a few-lines proof why there is no EPR Paradox within the framework of the New Relativity Theory if one assumes that information flows in a similar fashion as ordinary charges in Electromagnetism ; i.e information is to be thought of as a ” field ” . Interestingly enough, this will be our only assumption. We are not implying that there is such a thing as a ” fifth ” force in Nature found one morning in the closet of our homes after a bad night. We are just voicing out what has been irrefutably proven over and over by experiments.
Take an electron-positron pair colliding at the center O of an infinite dimensional sphere, $`S_D`$ for $`D\mathrm{}`$, at a givent moment we call $`t=0`$. After the collision a pair of two photons will travel in opposite directions imposed by energy-momentum conservation. An any given moment after the collision, we can locate those two photons at the surface of a multidimensional sphere of radius $`R=ct`$. The flux of information from the center of the sphere O flowing from the moment of the $`e^{}/e^+`$ collision radially outwards through the hypersurface is :
$$\mathrm{\Phi }=\stackrel{}{J}_D.d\stackrel{}{S}_{D1}=J_DS_{D1}.$$
$`(2)`$
This is nothing but the usual Gauss Law in Electromagnetism. The $`D`$-dimensional information-current, $`J_D`$, points radially outwards from the center O. Due to hyper-spherical symmetry its magnitude only depends on the radius $`R=ct`$. At each given point on the hypersurface, the current is pointing radially outwards and has the same value of magnitude, $`J_D(R)`$, along all the points of the hyper-sphere, This is why one can pull out the current outside the integral. The hypersurface $`S_{D1}`$ encloses inside a $`V_D`$ volume given in terms of gamma functions. Similar considerations apply to the higher-dimensional solid angle :
$$V_D=\frac{\pi ^{D/2}R^D}{\mathrm{\Gamma }(\frac{D+2}{2})}.S_{D1}=\frac{dV_D}{dR}=R^{D1}\mathrm{\Omega }_{D1}.\mathrm{\Omega }_{D1}=\frac{1}{R^{D1}}\frac{dV_D}{dR}.$$
$`(3)`$
Therefore, the total information-flux is given by the usual Gauss Law :
$$\mathrm{\Phi }=J_D(R).R^{D1}\mathrm{\Omega }_{D1}=J_D(R).R^{D1}.\frac{1}{R^{D1}}\frac{dV_D}{dR}=J_D(R).\frac{D\pi ^{D/2}R^{D1}}{\mathrm{\Gamma }(\frac{D+2}{2})}.$$
$`(4)`$
Now we take the $`D\mathrm{}`$ limit and make use of Stirling’s asymptotic formula for the gamma function : ( Pictures drawn on a Mathematica package also verify explicitly the results below )
$$lim_D\mathrm{}\mathrm{\Gamma }(\frac{D+2}{2})\sqrt{2\pi }(\frac{D+2}{2})^{\frac{D+2}{2}}e^{\frac{D+2}{2}}.$$
$`(5)`$
By Radius $`R=ct`$ one means radius in Planck scale units. We will set the Planck scale to 1. So by $`lnR`$ in all of the formulae below we mean $`ln(R/\mathrm{\Lambda })`$ Otherwise the units will not match up.
As $`D\mathrm{}`$ one can verify that in the asymptotic $`D=\mathrm{}`$ limit the numerator expression for the flux approaches :
$$exp[lnD+\frac{D}{2}ln\pi +(D1)lnR]exp[lnD+\frac{D}{2}ln\pi +DlnR]exp[lnD+Dln\pi +DlnR]$$
$`(6)`$
whereas the denominator approaches :
$$exp[\frac{D+2}{2}ln(\frac{D+2}{2})\frac{D+2}{2}]exp[DlnDD]exp[(D1)lnD]exp[DlnD].$$
$`(7)`$
Hence, the flux in the infinite $`D`$ limit is :
$$\mathrm{\Phi }=Jexp[lnD+Dln\pi +DlnRDlnD]=Je^\alpha .$$
$`(8a)`$
To be precise, upon reinserting the Planck scale one has that the flux is given in Planck units as :
$$\mathrm{\Phi }=J(\mathrm{\Lambda })^{D1}e^\alpha =Je^\alpha \times 1^{D1}=Je^\alpha .$$
$`(8b)`$
where $`\alpha `$ is :
$$\alpha =D(ln\pi +lnR)+(1D)lnDD[ln\pi +lnRlnD]$$
$`(9)`$
For $`finite`$ times , in units of $`\mathrm{\Lambda }=1`$, $`R=ct\mathrm{}`$ the coefficient $`\alpha `$ goes to negative infinity :
$$\alpha DlnD\mathrm{}e^\alpha 0$$
$`(10)`$
So
$$lim_D\mathrm{}J_D(R).\frac{D\pi ^{D/2}R^{D1}}{\mathrm{\Gamma }(\frac{D+2}{2})}J_{\mathrm{}}(R)\times 0=\mathrm{\Phi }J_{\mathrm{}}(R)\mathrm{}.$$
$`(11)`$
hence, in the $`D=\mathrm{}`$ limit, the current ( in Planck units, $`\mathrm{\Lambda }=1`$ ) blows up. This is not because there is a superluminal speed of information. It is because the hyper-volume, hyper-area elements, for finite values of $`R`$, go to zero in infinite dimensions!. Everything shrinks to a hyperpoint despite the fact that the radius is not zero ! The hyperpoint is the infinite-dimensional version of a point in ordinary finite-dimensional spacetime.
The current is as usual of the form : $`J=\rho v`$. As the hyper-volume, hyper-area elements, for finite values of $`R`$, go to zero, the information charge density $`\rho `$, charge per unit hyper-volume, blows up !. The information charge density diverges at the hyperpoint. The information velocity $`v`$ is constant and cannot exceed the speed of light. From the point of view of an infinite-dimensional observer, all the points of the hypersurface are interconnected. There is no such thing as non-locality in Quantum Mechanics. This is an illusion due to the shrinking to zero ( for finite radius) of the infinite-dimensional volume of the hypersphere, resulting from the asymptotic behaviour of the gamma functions !
This corroborates Mach’s brilliant insight that everything is connected in the ( Quantum ) Universe. What happens here and now, affects everything in the Universe in an instant. Based on the recent teletransportation experiments of a single photon by several experimental teams, this view of the Quantum Universe may lead an advance future generation of open minded scientists to achieve the ultimate communication system : instant exchange of information to anyplace in the Universe by tapping into the infinite dimensions of Quantum Spacetime. An speculative application of this would be to tele-transport a quantum copy of the human genome to other distant Planets in the Universe suitable for life. This would be a way out of the Galactic bounds we live in and an escape of the ultimate fate of the earth : consumed by the Sun when it becomes a Red Giant.
Spacetime travel in an instant will be much harder to achieve if by travel on means tele-transporting a quantum copy of ourselves to another point in the Universe. In order to do that one has to be able to tele-transport our consciousness as well. We adscribe to Penrose’s view that consciousness is a non-algorithmic process. This agrees with the Uncountably-infinite number of dimensions of the Cantorian-Fractal Spacetime view of El Naschie . It would be impossible for a Quantum Turing Machine ( a Quantum Computer) to quantum-process such vast of uncountably-infinite number of quantum bits. Never, in our wildest dreams we could possibly count such large number of dimensions of the Cantorian-Fractal Spacetime of El Naschie . Such World is not a mere Mathematical abstraction : it is essential for Consciousness to emerge. It is desirable that The Theory of ” Everyhthing ” should include Consciousnes. The Theory of Everything has to account for the existence of Conscious life and when, why, how, and for what it emerged from the Quantum Universe. A ”pointeless ” Universe is another one of those alarming signals that something is inherently incomplete with our view of the World. We believe that it is not sufficient to dismiss these questions as ”meaningless metaphyiscs ”.
Upon closer inspection of eq-(11), if one were to set $`J=finite`$ ; this would imply that the information flux $`\mathrm{\Phi }=0`$ so by Gauss Law there is $`nonet`$ information charge enclosed in the hypersphere. This is not correct for the following reason.
Nottale’s Scale Relativity implies that it is not possible to have zero measures with zero resolutions. It is possible to have zero measures but with ( nonzero) Planck scale resolutions. The $`e^{}/e^+`$ pair never goes beyond the minimum Plank scale resolution. The center O of the hypersphere is not a physical point. It is a smeared fuzzy hypersphere of infinite dimensions but with a nonzero Planck scale radius. This is a reason why Noncommutative Geometry, Fuzzy Phyiscs, Quantum Groups ….could be the right approaches to look at the world at small scales. Thus the information charge is distributed ” uniformly ” , in discrete bits of Planck hyper-area, in Planck units , over the outer ” surface ” of the hyperball of Planck radius. There is no inside. Inside is meaningless notion below the Planck scale, this is why the information charge has to reside on the ” surface ”. It would not be so surprising if this mechanism could be linked to the Bekenstein-Hawking entropy-area relationship of Black Holes.
The number of dimensions increases as one probes finer and finer $`resolution`$-scales ( not to be confused with lengths, although they both have the same units). By resolutions one means the resolutions that a physical apparatus can resolve. Resolutions which are not the same thing as the spacetime labels of a ” point ”, event ” like $`x^\mu `$. Resolutions that so far ( until Nottale) have been overlooked in the description of Physics. As one approaches asymptotically the Planck scale $`resolution`$ , the hypersphere of Planck scale radius becomes more and more ” visible” to us . To be able to reach this limiting ” threshold ” of $`resolutions`$ in our physical apparatus, an infinite amount Energy is required as Nottale has argued. By the same token that it takes an infinite amount of energy to accelerate a mass ( nonzero rest mass) from rest to the speed of light, it takes an infinite energy to probe Planck scale-resolutions. The final infinite-dimensional hypersphere, containing the information charge located at the ” origin ” O, shrinks to a hyperpoint of zero size , but finite Planck radius, The information charge density also diverges at the infinite-dimensional point : the hyperpoint of nonzero Planck radius. Exactly in the same way it did for hyperspheres of radius $`R=ct`$ upon taking the infinite dimension limit.
Concluding, the flux $`\mathrm{\Phi }`$ is not zero. There is a net information charge enclosed by the hypersphere with center O and radius $`R=ct`$.
Exactly the same argument occur if one asks the question : What does one of the photons ” see ” ? It will ” see ” the other photon at a distance $`l=2R=2ct`$ ( in Planck units) move away at the speed of light. Due to the Doppler effect, the frecuency will be redshifted so much that the photon will appear to be completely dark, with zero frecuency . For finite values of the radius, $`l=2R=2ct`$, the hypersphere centered at one of the photons will again shrink to a zero size , to a hyperpoint, in the infinite Dimensional ( large $`D`$ ) limit. Therefore, when a Macroscopic Observer with a Physical Apparatus measures the polarization of one particular photon, it will transfer its information to the other photon in an instant due to the fact that both of the photons have access to an extremely large number of dimensions in comparison to the macroscopic observers; i.e the photons truly live inside the hyperpoint. For this reason, they are able to exchange information in an instant without actually having a superluminal speed of information ! It is the information charge density ( and information current $`J`$) that diverges once again at the hyperpoint, and not the information speed. It is only an illusion due to the shrinking to zero of the hypersphere in the infinite dimensional limit.
Of course, the $`e^{}/e^+`$ pair does not attain such infinite energies to probe Planck scale resolutions, they come very close to each other but never reach the Planck scale. As they approach each other more and more dimensions become visible to them. Much more dimensions than the dimensions of the apparent one-dimensional world to a macroscopic observer looking at the line between the two emerging photons while performing his experiment. . Effectively, the number of dimensions of the world visible to the $`e^{}/e^+`$ pair, and the two emerging photons , is very high in comparison to the apparent $`D=1,D=4`$ of the macrospcopic observers , that for all practical purpopses, one can take the infinite dimensional limit of the gamma functions. The diagrams explicitly show that the hypervolume, hyperareas fall-off very rapidly to zero as $`D`$ moves far away from the $`D=4`$. It is not necessary to actually take the infinite dimensional limit too literate.
A related textbook issue is the following : Imagine a rapid moving observer passing by ourselves during the night while we are gazing at the stars. Due to the Lorentz contraction the celestial sphere that he experiences will naturally shrink with respect to us. It shrinks, but does it appear flattened ? The answer is no. One can view the Lorentz transformations in spinor terms as a $`SL(2,C)`$ Mobius transformation. Since the Mobius transformation maps circles to circles, the celestial sphere will have shrunk in radius only but it will not be flattened.
Similar analogy happens to the photon. What does a photon ” see ” ? Since we have said earlier that one cannot for certain answer such questions. We can only follow what we know so far : Due to the infinite Lorentz contraction the celestial sphere will shrink to a point. Doesn’t this contradict Nottale view that the Planck scale is the mimimal length ? The answer is no. Once again we have to take the variable dimensions of the Quantum World that a photon experiences. The photon is a quantum entity. Nobody can deny this. The photon of a given energy $`E=\mathrm{}\omega `$ will probe resolutions larger than the Planck scale. Rigorously speaking , we should write : $`E=\mathrm{}_{eff}(k^2)\omega `$. In we have shown that the New Relativity Theory demands an energy-dependent effective Planck constant so that $`[x,p]=i\mathrm{}_{eff}(k^2)`$ to reproduce the full blown Quantum Spacetime Uncertainty Relations that are more general than the String Uncertainty Relations : we have included the effects of all extended objects . It was shown rigorously why one cannot probe resolutions smaller than the Planck scale. As energy begins to be pump-in, one cannot probe smaller scales. Spacetime actually starts to grow. It is possible that a polymerization growth process of the Quantum Spacetime begins : an infinite chain of self similar branched polymers is triggered and baby universes branch off. The Quantum Universe might be an ever self-reproducing , self-recursive, self-iterated fractal process as Linde has suggested.
Only at infinite energy will a photon be able to probe the Planck scale. The celestial hypersphere that the photon ” sees ” has a radius of the order of the inverse photon Energy, roughly, assuming it is a low energy photon, Energy and resolution are inversely correlated at that level., not at higher enegy levels. Scale Relativity implies that the Compton wavelength and momentum are decoupled as one approaches Planck scales. It takes an infinite energy to probe the Planck scale. The Planck scale is the ultimate Ultaviolet Regulator. However, due to the effectively large number of dimensions that the photon has access to, despite the fact the hypersphere has a nonzero radius, the celestial hypersphere shrinks to zero size consistent with the infinite Lorentz contraction ! . It is true that one has to construct the full Scale-Motion Relativity to be fully rigourous and consistent. We have presented a solution to the apparent paradox of how one can have a zero measure/size ( due to the infinite Lorentz contraction) with a nonzero resolution for a radius : Infinite ( large number of ) Dimensions is the key once again !
Therefore, in essence : By introducing the notion of hyperpoint in physics, which is forced upon us by the New Relativity Principle as a result of having a truly infinite dimensional Quantum World. we have imbued a mathematical point with a true physical meaning : it is an infinite-dimensional hypersphere, of zero size but nonzero radius !.
When $`t=\mathrm{}`$ then the coefficient $`\alpha `$ will no longer be negative infinity due to a cancellation between $`lnR`$ and $`lnD`$ :
$$\alpha D[ln\pi +lnRlnD]=D[ln\pi +ln(ct)lnD]Dln\pi e^\alpha \mathrm{}.$$
$`(12)`$
In this case, one has the opposite result : the value of the information current $`J_D`$ at $`R=\mathrm{}`$ collapses to zero, as it should on physical grounds. The information field must vanish at infinity in any dimension, finite or infinite. As the photons move away from eachother, if one waits an infinite amount of time to peform the EPR experiment, the photons will no longer be correlated !
To sum up :
The EPR Paradox only occurs to the one-dimensional beings ( or finite-dimensional beings) living along the linear path ( around the linear path) of the photons who wish to perform the EPR gedanken experiment. From their finite-dimensional point of a view, QM appears to be non-local : a superluminal transfer of information appears to take place. From the point of view of the New Relativity Theory there is no paradox because Quantum Spacetime is truly infinite-dimensional. For those Quantum-dimensional beings who were able to tap into the effectively ”infinite” number of dimensions of Quantum Spacetime at the very ” moments ” when the $`e^{}/e^+`$ pair collided, at a very small distance separation among them, distance which cannot be smaller than the Planck scale as indicated by Nottale’s Scale Relativity, there is no such Paradox at all ! : the information current blows up because from their infinite-dimensional point of view, for finite values of the radius, the hypersphere has shrunk to a hyperpoint. The transfer of information to the two photons, about the spin and other quantum numbers of the $`e^{}/e^+`$ pair, occurs in an instant ! Every point in their universe is inter-connected as Mach argued long ago. Similar arguments apply to the two photons when a macroscopic observer measures the polarization of one photon, the information is transfered to the other photon in an instant via an effectively ” infinite ” dimensional ( relative to the macroscopic observers) Quantum Spacetime accesible to them.
This should encourage us to view Feynman’s path integral formulation of QM taking all posible paths in a finite dimensional spacetime, from the New Relativity Theory point of view : it is possible to have a finite number of paths in an Infinite Dimensional Quantum Spacetime. The main question is : Where does the Feynman statistical complex-weighting of the paths via the $`e^{iS}`$ comes from ?
The partial answer was given by Ord , Nottale and others :
Since fractal paths have a dominant weight in the path integral compared to the smooth ones, the latter have a zero measure compare to the former, roughly speaking, Quantum Effects manifest or channel themselves via the fractality of spacetime. Although there are people who do not subscribe to this view. Fractal curves are continous but nowhere differentiable. This means that the derivatives are discrete-valued. The discrete jumps of the values of the tangents are ”quantized” in units of what has been called by mathematicians the ” Planck ” constant of a curve. In this fashion the Feynamn $`e^{iS}`$ weighting factor is interpreted although , we must say that no rigorous proof of this has been given as far as we know.
Fractals and Scale Relativity are essential because as the resolutions that a physical apparatus can resolve reach the mimimal Planck scale resolution ( resolutions must not to be confused with statistical uncertainties nor with ordinary lengths) the number of fractal dimensions blows up. For a new Phase space path integral derivation of Feynman’s particle propagator that is $`roughly`$ based on these ideas that a fractal particle ” path ” can have a meaning in QM see .
The apparent superluminal information velocity happens in other aspects of Physics. There is a very simple analogy with superluminal jets in Astrophysics . If one takes a flash light at a sufficiently large distance from a wall and rotates it very rapidly , the image on the wall can appear to move faster than light. However the image is not a truly physical object. The physical photons never move faster than light. The image is comprised of many different photons and not of a fixed particular number of them . The maximum angular velocity of rotation of the flash light is bounded by Special Relativity :
$$\omega _{max}=\frac{c}{r}$$
$`(13)`$
where $`r`$ is the length of the flashlight. If the distance to the wall is $`R`$ then the apparent velocity of the shadow is :
$$v=\frac{cR}{r}>c.$$
$`(14)`$
The ( unphysical object) shadow can move faster than light. One does not even have to go to such extremes of achieving the maximal angular speed for the flash light , one can simply choose the wall far enough, and the flashlight sufficently bright, ( $`R`$ large enough ) so the image on the wall moves with a superluminal velocity $`\omega R>c`$. Exactly similar arguments occur with the phase velocity in wave propagation. The phase velocity can be greater that $`c`$ but the physical group velocity is always bounded by $`c`$. Taking the number of Dimensions to infinity, mimics this simple example of taking the distance to the wall far enough and rotating the flash light fast enough.
Similar arguments can be taken with the so called Back Hole Information loss Paradox. Since Quantum Spacetime is truly infinite dimensional, there is no such thing as an Information Loss. This information is stored in all the infinite number of dimensions that are inaccesible to an outside low energy observer. There is information radiated away and a remnant ” hidden ” in the infinite number of dimensions inaccesible to the outside observers. Black Hole evaporation stops at the minimal scale in Nature : the Planck scale, reaching a maximum temperature, Planck’s Temperature. Scale Realtivity not only induces an effective value of the Planck’s constant : $`\mathrm{}_{eff}(k^2)`$ , it also affects the Boltzmann constant as well : $`k_B(k^2)`$ so that : $`k_B(k^2)T=\mathrm{}_{eff}(k^2)\omega `$. As one reaches the Planck scale, energy blows up but the temperature reaches asymptotically the maximum Planck Temperature ( thermal Relativity). One must have a standard of temperature to compare temperatures with. That maximum universal standard is the Planck Temperature whose definition in $`D=4`$ is :
$$T_P=\sqrt{\frac{\mathrm{}c^5}{Gk_B^2}}=1.42\times 10^{32}K.$$
$`(15)`$
Astrophysicits have been baffled by recent findings that there are unexplained extremely bright and unrelenting sources of energy. It is warranted to study these phenomena within the framework of the New Relativity Theory. To be able to ” see ” all the information one has to tap into all the infinite number of dimensions of the Quantum Spacetime. To achieve that one requires infinite amount of energy to probe the Planck scale resolutions according to the Scale Relativity Principle . At that scale (infinite) Dimensions, (infinite) Energy and (infinite) Information merges into the ” Omega ” hyperpoint , the ” Trinity ” hyperpoint….. the ultimate infinite-dimensional point : At that scale, the ” Trinity ” hyperpoint, Dimensions, Energy and Information are indistinguishable from each other. More details will be given later.
Acknowledgements
We are indebted to E. Spallucci for a very constructive critical remarks. We thank G. Chapline. L.Nottale , W. Pezzaglia, M. El Naschie and D. Finkelstein for illuminating discussions. Finally many thanks to C. Handy and M. Handy for their assistance and encouragement .
References
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hep-th/0001023. ”
” Hints of a New Relativity Principle from $`p`$-brane Quantum Mechanics ” hep-th/9912113.
” Is Quantum Spacetime Infinite Dimensional ? hep-th/0001134.
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in the $`^{(\mathrm{})}`$ Space. Jour. Chaos. Solitons and Fractals 11 (2000) 1149-1162. |
warning/0002/hep-ph0002177.html | ar5iv | text | # I INTRODUCTION
## I INTRODUCTION
A major theoretical challenge posed by the solutions to the atmospheric and solar neutrino anomalies is that the atmospheric neutrino data require a large $`\nu _\mu \nu _\tau `$ mixing, whereas the corresponding quark mixing between the second and the third generation is very small. This is not easy to understand in the context of quark-lepton unified theories. While there are suggestions to understand a large mixing in the context of various kinds of unified theories including the SO(10), where there is a natural quark-lepton unification, no convincing natural model has yet emerged. It is therefore necessary to explore alternative possibilities. One way to proceed is not to concentrate on a particular model, but to look for the features that a model should have in order to be able to predict the observed large mixing naturally.
In a recent paper , we pointed out that for two Majorana neutrinos with the same CP parity that are nearly degenerate in mass, a small neutrino mixing at the high scale can be magnified by the radiative corrections through the renormalization group running down to the weak scale. In such theories, there would be no need to put special constraints on the mixings in the theory at the high (e.g. seesaw) scale and indeed the quark and lepton mixings can be very similar (as, say, would be predicted by the simple seesaw models). It is the goal of this paper to show that such a mixing pattern, which involves only small mixings at the high scale ($`\mathrm{\Lambda }`$), can explain the neutrino anomalies at the low scale as long as the conditions outlined in are satisfied. It is then possible to explain the solar neutrino problem through the small angle MSW solution and the atmospheric neutrino problem through the large mixing angle, which gets generated through the radiative magnification.
The paper is organized as follows. In Sec. II, we introduce the $`(\mathrm{\Omega },\mathrm{\Phi },\mathrm{\Psi })`$ parametrization for the mixing angles, and taking $`\mathrm{\Phi }=0`$ at the high scale, show that the radiative corrections can magnify $`\mathrm{\Psi }`$ while keeping $`\mathrm{\Omega }`$ and $`\mathrm{\Phi }`$ unaffected. In Sec. III, we show that the condition $`\mathrm{\Phi }=0`$ is consistent with the current data from the solar, atmospheric and reactor experiments. In Sec. IV, we consider the possible 4$`\nu `$ mixing schemes that can explain the LSND results in addition, and identify the scenarios for which radiative magnification can provide a natural explanation through the quark-lepton unified theories. Sec. V concludes.
## II Radiative magnification for three neutrino mixing
In the absence of $`CP`$ violation in the lepton sector, the mixing matrix $`U_\mathrm{\Lambda }`$ at the scale $`\mathrm{\Lambda }`$ can be parametrized as
$$U_\mathrm{\Lambda }=U_{12}(\mathrm{\Omega })\times U_{13}(\mathrm{\Phi })\times U_{23}(\mathrm{\Psi }),$$
(1)
where all the three rotation angles lie between 0 and $`\pi /2`$. Note that the order of multiplication of the rotation matrices is different from the conventional one , so the angles $`\mathrm{\Omega },\mathrm{\Phi },\mathrm{\Psi }`$ involved here should not be mistaken for the angles $`\omega ,\varphi ,\psi `$ used conventionally. Nevertheless, (1) is a perfectly valid way of parametrizing the mixing matrix, and is useful for addressing a certain class of problems (e.g. see ).
At the low scale $`\mu `$, the mixing matrix $`U_\mu `$ can be written in general as
$$U_\mu =U_{12}(\overline{\mathrm{\Omega }})\times U_{13}(\overline{\mathrm{\Phi }})\times U_{23}(\overline{\mathrm{\Psi }}).$$
(2)
The CHOOZ results indicate a small $`U_{e3}`$, which corresponds to a small value for
$$\mathrm{cos}\overline{\mathrm{\Omega }}\mathrm{cos}\overline{\mathrm{\Psi }}\mathrm{sin}\overline{\mathrm{\Phi }}+\mathrm{sin}\overline{\mathrm{\Omega }}\mathrm{sin}\overline{\mathrm{\Psi }}.$$
This can be satisfied with the choice of $`\overline{\mathrm{\Phi }}=0`$ and a small $`\mathrm{sin}\overline{\mathrm{\Omega }}\mathrm{sin}\overline{\mathrm{\Psi }}`$. That such a choice can satisfy the solar and the atmospheric data is shown in Sec. III. With this motivation, we start with $`\mathrm{\Phi }=0`$ at the high scale (this choice leads to $`\overline{\mathrm{\Phi }}=0`$, as we shall show in this section), and show that the radiative corrections can magnify $`\mathrm{\Psi }`$ while keeping $`\mathrm{\Omega }`$ and $`\mathrm{\Phi }`$ unaffected.
With only the $`\mathrm{\Omega }`$ and $`\mathrm{\Psi }`$ mixings nonzero at the scale $`\mathrm{\Lambda }`$, the effective mass matrix $`M_\mathrm{\Lambda }^{eff}`$ in the flavor basis is
$$M_\mathrm{\Lambda }^{eff}=U_\mathrm{\Lambda }M_\mathrm{\Lambda }^dU_\mathrm{\Lambda }^{}=U_{12}(\mathrm{\Omega })U_{23}(\mathrm{\Psi })M_\mathrm{\Lambda }^dU_{23}^{}(\mathrm{\Psi })U_{12}^{}(\mathrm{\Omega }),$$
(3)
where $`M_\mathrm{\Lambda }^d=Diag(m_1,m_2,m_3)`$. If the radiative corrections are included , we have
$$M_\mathrm{\Lambda }^{eff}M_\mu ^{eff}=\left(\begin{array}{ccc}\sqrt{I_e}& 0& 0\\ 0& \sqrt{I_\mu }& 0\\ 0& 0& \sqrt{I_\tau }\end{array}\right)M_\mathrm{\Lambda }^{eff}\left(\begin{array}{ccc}\sqrt{I_e}& 0& 0\\ 0& \sqrt{I_\mu }& 0\\ 0& 0& \sqrt{I_\tau }\end{array}\right),$$
(4)
where $`I_\alpha 12\delta _\alpha `$ are the radiative corrections that appear due to the Yukawa couplings of the charged leptons $`e,\mu `$ and $`\tau `$ respectively. Given the strong hierarchical pattern of the charged lepton masses, we neglect the corrections due to $`e`$ and $`\mu `$, i.e. $`I_e=I_\mu =1`$. Let us define $`_\tau Diag(1,1,\sqrt{I_\tau })`$. Then from (3) and (4),
$$M_\mu ^{eff}=_\tau U_{12}(\mathrm{\Omega })U_{23}(\mathrm{\Psi })M_\mathrm{\Lambda }^dU_{23}^{}(\mathrm{\Psi })U_{12}^{}(\mathrm{\Omega })_\tau .$$
(5)
Noting that $`[U_{12}(\mathrm{\Omega }),_\tau ]=0`$, we get
$$M_\mu ^{eff}=U_{12}(\mathrm{\Omega })[_\tau U_{23}(\mathrm{\Psi })M_\mathrm{\Lambda }^dU_{23}^{}(\mathrm{\Psi })_\tau ]U_{12}^{}(\mathrm{\Omega }).$$
(6)
The quantity in the square brackets in (6) is in a form where the first row and column are effectively decoupled and the situation reduces to the two-generation mixing, which has been considered in detail in . This quantity can be written as
$$_\tau U_{23}(\mathrm{\Psi })M_\mathrm{\Lambda }^dU_{23}^{}(\mathrm{\Psi })_\tau =U_{23}(\stackrel{~}{\mathrm{\Psi }})M_\mu ^dU_{23}^{}(\stackrel{~}{\mathrm{\Psi }}),$$
(7)
where $`M_\mu ^d`$ is a diagonal matrix. The new (2-3) mixing angle $`\stackrel{~}{\mathrm{\Psi }}`$ is given by
$`\mathrm{tan}(2\stackrel{~}{\mathrm{\Psi }})`$ $`=`$ $`{\displaystyle \frac{\mathrm{tan}(2\mathrm{\Psi })}{\lambda }}(1+\delta \tau ),`$ (8)
$`\lambda `$ $``$ $`{\displaystyle \frac{(m_3m_2)C_{2\mathrm{\Psi }}+2\delta _\tau m}{(m_3m_2)C_{2\mathrm{\Psi }}}},`$ (9)
where $`m`$ is the common mass of the quasi-degenerate neutrinos. Now, if
$$\delta _\tau \frac{(m_2m_3)C_{2\mathrm{\Psi }}}{2m},$$
(10)
then $`\lambda 0`$, so that the mixing angle $`\stackrel{~}{\mathrm{\Psi }}`$ becomes large . Since $`\delta _\tau 1`$, for the condition (10) to be satisfied, $`\nu _2`$ and $`\nu _3`$ need to have the same CP parity. Thus the $`\mathrm{\Psi }`$-mixing can be magnified at the weak scale, which explains the atmospheric neutrino data (See Sec. III).
From (6) and (7),
$$M_\mu ^{eff}=U_{12}(\mathrm{\Omega })U_{23}(\stackrel{~}{\mathrm{\Psi }})M_\mu ^dU_{23}^{}(\stackrel{~}{\mathrm{\Psi }})U_{12}^{}(\mathrm{\Omega }).$$
(11)
This shows that the same (1-2) mixing angle $`\mathrm{\Omega }`$ that was needed for diagonalizing $`M_\mathrm{\Lambda }^{eff}`$ is also needed for diagonalizing $`M_\mu ^{eff}`$ \[see (3) and (11)\], and that a (1-3) mixing angle $`\mathrm{\Phi }`$ is not required. Thus, $`\overline{\mathrm{\Psi }}=\stackrel{~}{\mathrm{\Psi }}`$, $`\overline{\mathrm{\Omega }}=\mathrm{\Omega }`$ and $`\overline{\mathrm{\Phi }}=\mathrm{\Phi }=0`$ are the mixing angles at the low scale.
As we shall see in Sec. III, we can explain the solar, atmospheric and the CHOOZ data with $`\overline{\mathrm{\Psi }}\pi /4`$ and a small $`\overline{\mathrm{\Omega }}`$ (corresponding to the SMA solution for the solar neutrinos). In a typical quark-lepton unified theory, $`\mathrm{\Omega }`$ would be small at the high scale. In the limit of neglecting the radiative corrections due to the second generation (i.e. $`I_\mu 1`$) that we have considered here, the magnification of $`\mathrm{\Omega }`$ due to radiative corrections is not possible. Also, if the $`CP`$ parity of the neutrino $`\nu _1`$ is opposite to that of $`\nu _2`$ and $`\nu _3`$ (which is required to ascertain the stability of a possible small nonzero $`\mathrm{\Phi }`$), a small $`\mathrm{\Omega }`$ at the high scale will stay small even when the radiative corrections due to $`\mu `$ are taken into account. Thus, the stability of a small $`\mathrm{\Omega }`$ is guaranteed, and the small angle MSW scenario can be generated naturally within the unification models.
The radiative corrections from the second generation \[i.e. $`_\mu Diag(1,\sqrt{I_\mu },1)Diag(1,1,1)`$\] modify (6) to
$$M_\mu ^{eff}=_\mu U_{12}(\mathrm{\Omega })[_\tau U_{23}(\mathrm{\Psi })M_\mathrm{\Lambda }^dU_{23}^{}(\mathrm{\Psi })_\tau ]U_{12}^{}(\mathrm{\Omega })_\mu .$$
(12)
Since $`[U_{12}(\mathrm{\Omega }),_\mu ]0`$, the value of $`\mathrm{\Omega }`$ may now get modified and $`\mathrm{\Phi }`$ may get generated. The value of $`\overline{\mathrm{\Psi }}`$ is also different from the value of $`\stackrel{~}{\mathrm{\Psi }}`$ as given in (9). But since
$$[U_{12}(\mathrm{\Omega }),_\mu ]=\delta _\mu \mathrm{sin}\mathrm{\Omega }\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right),$$
(13)
and the values of $`\mathrm{\Omega }`$ and $`\delta _\mu `$ are both small, these differences $`(\overline{\mathrm{\Omega }}\mathrm{\Omega })`$, $`(\overline{\mathrm{\Phi }}\mathrm{\Phi })`$ and $`(\overline{\mathrm{\Psi }}\stackrel{~}{\mathrm{\Psi }})`$ are not expected to be large.
## III Satisfying the solar, atmospheric and CHOOZ data
In the following, we show that our choice of parametrization (1) with $`\mathrm{\Phi }=0`$ can explain the solar and atmospheric anomalies and still be consistent with the stringent bounds coming from the CHOOZ experiment.
We first concentrate on the CHOOZ and the atmospheric data which share a common mass scale $`\mathrm{\Delta }m_{31}^2\mathrm{\Delta }m_{32}^210^3eV^2`$. In this case, the relevant probability expressions are
$`P_{\alpha \beta }^{atm}`$ $``$ $`4|U_{\alpha 3}|^2|U_{\beta 3}|^2\mathrm{sin}^2\left({\displaystyle \frac{\mathrm{\Delta }m_{31}^2L}{4E}}\right),`$ (14)
$`P_{\alpha \alpha }^{atm}`$ $``$ $`14|U_{\alpha 3}|^2(1|U_{\alpha 3}|^2)\mathrm{sin}^2\left({\displaystyle \frac{\mathrm{\Delta }m_{31}^2L}{4E}}\right).`$ (15)
To satisfy the CHOOZ constraint
$$|U_{e3}|^2<0.03\text{ for }\mathrm{\Delta }m_{31}^2>210^3\text{ eV}^2$$
(16)
with $`\overline{\mathrm{\Phi }}=0`$, we need
$$\mathrm{sin}\overline{\mathrm{\Omega }}\mathrm{sin}\overline{\mathrm{\Psi }}<0.17.$$
(17)
This small value of $`|U_{e3}|^2`$ also guarantees $`P_{ee}1`$ \[eq. (15)\] and $`P_{\mu e}0`$ \[eq. (14)\] in the atmospheric neutrino data.
A fit to the $`L/E`$ distribution of the atmospheric neutrinos gives at 90% confidence level, using (15),
$$0.2<|U_{\mu 3}|^2<0.8.$$
(18)
This corresponds to
$$0.45<\mathrm{cos}\overline{\mathrm{\Omega }}\mathrm{sin}\overline{\mathrm{\Psi }}<0.9.$$
(19)
By examing the mixing matrix as parametrized in (2), we can see that $`\overline{\mathrm{\Phi }}=0`$, $`\overline{\mathrm{\Psi }}\pi /4`$ and a small $`\overline{\mathrm{\Omega }}`$ can easily satisfy the requirements (17) and (19). The smallness of $`\overline{\mathrm{\Omega }}`$ is forced by the CHOOZ constraints and the large value of $`\mathrm{sin}\overline{\mathrm{\Psi }}`$ is required by the atmospheric $`P_{\mu \mu }`$.
Let us now consider the solar neutrino solution. In the case of the solar neutrino anomaly, the SMA solution corresponds to
$$|U_{e2}|^2(0.5÷2.5)10^3,$$
(20)
whereas the other solutions – LMA, LOW and VO – correspond to $`|U_{e2}|^2>0.2`$. In the parametrization (2),
$$|U_{e2}|^2=\mathrm{sin}^2\overline{\mathrm{\Omega }}\mathrm{cos}^2\overline{\mathrm{\Psi }}.$$
(21)
It is difficult to reconcile the smallness of $`\overline{\mathrm{\Omega }}`$ forced by the atmospheric and CHOOZ results to the $`|U_{e2}|^2`$ required for the LMA, LOW or VO solution. But in the case of the SMA solution, (20) and (21) give
$$\mathrm{sin}\overline{\mathrm{\Omega }}\mathrm{cos}\overline{\mathrm{\Psi }}0.02÷0.05,$$
(22)
which can be satisfied simultaneously with (17) and (19). The region in the $`\overline{\mathrm{\Omega }}\overline{\mathrm{\Psi }}`$ parameter space that satisfies all the constraints (17), (19) and (22) is shown in Fig. 1. Our scheme thus supports the SMA solution: if we start with a small $`\mathrm{\Omega }`$ at the high scale (which is natural in the quark-lepton unified theories), it does not change much through radiative corrections (as we have shown in sec. II), and a small $`\overline{\mathrm{\Omega }}`$ is retained at the low scale.
As pointed out in , the mechanism of radiative magnification does not need any fine-tuning, but is at work in a range of parameter space for any given model. As an example of the radiative magnification of $`\mathrm{\Psi }`$, let us consider MSSM, where the parameter $`\mathrm{tan}\beta `$ determines the magnitude of the radiative corrections. The value of $`\mathrm{tan}\beta `$ required to obtain any given magnified value of $`\overline{\mathrm{\Psi }}`$ is shown in Fig. 2. This indicates the phenomenologically interesting range of $`\mathrm{tan}\beta `$ for radiative magnification.
## IV Four neutrino schemes
The features of radiative magnification noted here can be used in order to identify the 4$`\nu `$ mixing scenarios in which the large atmospheric mixing can be naturally generated. Taking into account that the recent atmospheric neutrino results disfavor ($`\nu _\mu \nu _s`$) oscillations , the 4$`\nu `$ solution for all the anomalies (atmospheric , solar and LSND ) is essentially of the form
$$[\nu _e\nu _s][\nu _\mu \nu _\tau ],$$
(23)
where the $`[\nu _e\nu _s]`$ pair ($`\mathrm{\Delta }m_{14}^2\mathrm{\Delta }m_{}^2`$) and the $`[\nu _\mu \nu _\tau ]`$ pair ($`\mathrm{\Delta }m_{23}^2\mathrm{\Delta }m_{atm}^2`$) are separated by $`\mathrm{\Delta }m_{es\mu \tau }^2\mathrm{\Delta }m_{LSND}^2`$. The solar neutrino puzzle is explained by the $`\nu _e\nu _s`$ oscillations and the atmospheric data are explained by the $`\nu _\mu \nu _\tau `$ oscillations. A small $`\nu _e\nu _\mu `$ mixing then explains the LSND observations.
In (23), the neutrinos can be considered to be written in the increasing order of masses. With the current data, it is still possible to change the order of neutrinos within a bracket, or the order of the brackets themselves. The order within a bracket will not have any influence on our conclusions, so we have only the two independent cases: (a) $`m_{es}<m_{\mu \tau }`$ and (b) $`m_{es}>m_{\mu \tau }`$, where $`m_{es}`$ ($`m_{\mu \tau }`$) denotes the average mass of the $`[\nu _e\nu _s]`$ ($`[\nu _\mu \nu _\tau ]`$) pair. In the case (a), $`\nu _\mu `$ and $`\nu _\tau `$ are necessarily quasi-degenerate: taking $`\mathrm{\Delta }m_{LSND}^21`$ eV<sup>2</sup> and $`\mathrm{\Delta }m_{atm}^24\times 10^3`$ eV<sup>2</sup>, we get the degree of degeneracy ($`\frac{\delta m}{m}`$) for the $`\nu _\mu \nu _\tau `$ pair as $`\frac{\delta m}{m}<2\times 10^3`$. Then the $`\mu \tau `$ mixing angle $`\theta _{\mu \tau }`$ can be radiatively magnified, as we require for the atmospheric neutrino solution.
In the case (b), the neutrinos $`\nu _\mu `$ and $`\nu _\tau `$ need not be quasi-degenerate, so the magnitude of radiative corrections needed to magnify $`\theta _{\mu \tau }`$ is large. Accounting for the large $`\theta _{\mu \tau }`$ through radiative magnification is then difficult. Thus, if radiative magnification is the reason for the large $`\theta _{\mu \tau }`$, then the case (a) is favored, i.e. $`m_{es}<m_{\mu \tau }`$ on the grounds of naturalness.
## V Conclusions
We have shown that, with the parametrization $`U=U_{12}(\mathrm{\Omega })\times U_{13}(\mathrm{\Phi })\times U_{23}(\mathrm{\Psi })`$ of the lepton mixing matrix, $`\overline{\mathrm{\Phi }}=0`$, $`\overline{\mathrm{\Psi }}\pi /4`$ and a small $`\overline{\mathrm{\Omega }}`$ at the low scale can satisfy all the constraints from the solar, atmospheric and CHOOZ data. These mixing angles can be generated at the high scale with $`\mathrm{\Phi }=0`$ and small $`\mathrm{\Omega }`$ and $`\mathrm{\Psi }`$ (which is natural in the quark-lepton unified theories), and magnifying $`\mathrm{\Psi }`$ through radiative corrections while keeping $`\mathrm{\Omega }`$ and $`\mathrm{\Phi }`$ unaffected.
Let us add a few words on the realization of this scenario of radiative magnification in the unified theories. We have not given any specific model realization, rather we have pointed out a class of models that would be successful in generating a large lepton mixing naturally, starting from a small mixing at the high scale. It is not hard to see that such small mixing angle patterns can emerge at the high scale $`\mathrm{\Lambda }`$ in quark-lepton unified theories of type $`SU(2)_L\times SU(2)_R\times SU(4)_c`$ if the right-handed neutrino coupling is assumed to be an identity matrix since the Dirac mass matrix for neutrinos that goes into the seesaw matrix is then identical to the up-quark mass matrix. Thus even though our discussion in this paper is completely model independent, its realization in the context of unified theories is quite straightforward. Our work thus demonstrates a way to have a natural solution for the neutrino anomalies in the quark-lepton unified theories.
Acknowledgements
We thank WHEPP-6, Chennai, India, where this work was initiated. The work of RNM is supported by the NSF Grant no. PHY-9802551. The work of MKP is supported by the project No. 98/37/9/BRNS-cell/731 of the Govt. of India. A. D. would like to thank A. De Gouvea for helpful discussions. Balaji wishes to thank S. Uma Sankar for useful clarifications. |
warning/0002/quant-ph0002019.html | ar5iv | text | # PHYSICAL MODEL OF DIRAC ELECTRON. CALCULATION OF ITS MASS AT REST AND OWN ELECTRIC AND MAGNETIC INTENSITIES ON ITS MOMENT LOCATION.
## 1 Introduction.
The successful scientific research of some natural phenomenon is very often connected with some necessary idealization and soma minimum simplification of the phenomenon under investigation. Many marvellous phenomena and remarkable properties of the sunstance have been described by help of the powerful logic of the Quantum Theory (QngThr). The physical model (PhsMdl) of some physical phenomenon presents as an actual ingradient of the physical theory (PhsThr). This is a scientific way for construction of the (PhsMdl) of the Dirac electron (DrEl). Although till now nobody knows what the elementary particle (ElmPrt) means, there exist a possibility for a obvious consideration of the unusual behaviour of all the relativistic quantized micro particle (QntMcrPrt) by means of our transparent surveyed PhsMdl of the DrEl. The PhsMdl of the DrEl is offered in all my work in resent nineteen for bring of light to physical cause of the uncommon relativistic quantum behaviour of the DrEl and give the thru physical interpretation and sense of its dynamical parameters. It turns up that all the relativistic dynamical properties of the DrEl are results of the participate of its fine spread (FnSpr) elementary electric charge (ElmElcChrg) in the Schrodinger’s self-consistent fermion strong correlated vortical harmonic oscillations, then all the quantized dynamical properties of the SchEl are results of the participate of its well spread (WllSpr) ElmElcChrg in the nonrelativistic Furthian quantized stochastic boson circular harmonic oscillations. It is used as for a obvious teaching the occurred physical micro processes within the investigated phenomena, so for doing them equal with the capacity of its mathematical correct description by the mathematical apparatus of the both the quantum mechanics (QntMch): the nonrelativistic (NrlQntMch) and relativistic (RltQntMch).
The object of this lecture is to discuss, explain and bring to light on the physical interpretation of the nonrelativistic quantum behaviour of the Schrodinger’s electron (SchEl) and of the relativistic quantum behaviour of the dynamical parameters of the Dirac’s electron (DrEl). The PhsMdl of the DrEl is proposed by me twenty years ago. This PhsMdl can equally explain as the physical causes for its unusual classical stochastic and so the quantum dualistic wave-corpuscular behaviour of SchEl. One gives a new cleared picturesque physical interpretation with mother wit of the physical scene of the relativistic dynamical parameters of DrEl. In our transparent surveyed PhsMdl of the DrEl one will be regarded as some point like (PntLk) ElmElcChrg, taking simultaneously part in the following four different motions: A) The isotropic three-dimensional relativistic quantized (IstThrDmnRltQnt) Einstein stochastic boson harmonic shudders (EinStchBsnHrmShds) as a result of momentum recoils (impulse kicks), forcing the charged QntMcrPrt at its continuous stochastical emissions and absorptions of own high energy (HghEnr) virtual photons (VrtPhtns) by its PntLk ElmElcChrg). This jerky motion display almost Brownian classical stochastic behaviour (BrnClsStchBhv) during a small time interval $`\tau _1`$, much less then the period $`T`$ of the IstThrDmnRltQnt EinStchBsnHrmShds and more larger then the time interval $`t`$ of the stochastically emission or absorption of the Hgh-Enr VrtPhtn by its PntLk ElmElcChrg. In a consequence of such jerks along the IstThrDmnRltQnt EinStchBznHrmShds ”trajectory” the DrEl’s PntLk ElmElbChrg takes form of the fine spread (FnSpr) ElmElcChrg. B) The IstThrDmnRltQnt Schrodinger fermion vortical harmonic oscillation motion (SchFrmVrtHrmOscMtn) of the DrEl’s FnSpr ElmElcChrg. In a consequence of such jerks along the EinStchHrmOscMtn ”trajectory” the ”trajectory” of the DrEl’s FnSpr ElmElcChrg, participating in the IstThrDmnRltQnt SchFrmVrtHrmOscMtn takes a strongly broken shape. Only after the correspondent averaging over the ”trajectory” of the IstThrDmnRltQnt EinStchBsnHrmShdMtn we may obtain the cylindrically spread ”trajectory” of the IstThrDmnRltQnt SchFrmVrcHrmOscs’ one, having got the form of the crooked figure of an eight. Only such a motion along a spread uncommon ”trajectory” of the DrEl’s FnSpr ElmElcChrg could through a new light over the SchEl’s WllSpr ElmElcChrg’s space distribution and over the spherical symmetry of the SchEl’s WllSpr ElmElcChrg. This self-consistent strongly correlated IstThrDmnRltQnt SchFrmVrtHrmOscs’ motion may be described correctly by means of the four components of its total wave function (TtlWvFnc) $`\mathrm{\Psi }`$ and four Dirac’s matrices ; $`\alpha _j(\gamma _j)`$ and $`\beta (\gamma _o)`$.
## 2 Description of the relativistic quantized behaviour of the DrEl.
In this way we may do as well as make a possibility for making clear the spinor character of such a movement and all its consequences as the proper mechanical momentum (MchMmn)(spin) and the rest self-energy, fermion symmetry and fermion statistics. Only in a result of the participating the FnSpr ElmElcChrg in the IstThrDmnRltQnt SchFrmVrtHrmOscs’ motion all the components of the resultant self-consistent (RslSlfCns) ElcInt of own QntElcMgnFld may be exactly compensated and have zero values and all the components of the RslSlfCns MgnInt of own QntElcMgnFld would be doubled in a comparison with the corresponding RslSlfCns values of the MgnInt of own ClsElcMgnFld of the NtnClsMcrPrt with same WllSpr ElmElcChrg, fulfilling Einstein’s relativistic stochastic harmonic oscillations motion in conformity with the laws of the Einstein’s relativistic classical mechanics (RltClsMch). It is turn out that in a result of own useful participation of the DrEl’s FnSpr ElmElcChrg the self-energy at a rest of its electromagnetic self-action (ElcMgnSlfAct) between its continuously moving potential and vector-potential to be minimized in the with own corresponding self-consistent way. It turns up that the RslSlfCns values of the ElcInt and the MgnInt of own QntElcMgnFld of the DrEl are generated by its FnSpr ElmElcChrg, editing incessantly Hgh-Enr StchVrtPhts at different moments of the recent half period of time in its corresponding ivarious spatial positions and absorbed it in the form of the ElcMgnSlfAct in the point of the DrEl’s FnSpr ElmElcChrg’s instantaneous positions; C) The isotropic three-dimensional nonrelativistic quantized (IstThrDmnNrlQnt) Furthian stochastic boson vortical harmonic oscillations (FrthStchBsnVrtHrmOscs) of the SchEl as a result of the permanent electric interaction (ElcIntAct) of its WllSpr ElmElcChrg with the ElcItn of the resultant QntElcMgnFld of the low energy (LwEnr) StchVrtPhtns, stochastically generated by dint of the fluctuating vacuum (FlcVcm) in the form of exchanging StchVrtPhtns between FlcVcm and its FnSpr ElmElcChrg. This Furthian quantized stochastic uncommon behaviour of the SchEl with own participation in the random trembling motion (RndTrmMtn) is very similar to Brownian classical stochastic behaviour of the BrnClsMcrPrt with own participation in the RndTrmMtn. But in principle the exact description of the resultant behaviour of the SchEl owing of its participation in the IstThrDmnNrlQnt FrthStchBsnCrcHrmOsc motions could be done only by means of the NrlQntMch’s and nonrelativistic ClsElcDnm’s laws. D) The classical motion of the Lorentz’s electron (LrEl) around an well contoured smooth thin trajectory which is realized in a consequence of some classical interaction (ClsIntAct) of its overspread (OvrSpr) ElmElcChrg, bare mass or magnetic dipole moment (MgnDplMnt) with the intensity of some external classical fields (ClsFlds) as in the Newton nonrelativistic classical mechanics (NrlClsMch) and in the nonrelativistic classical electrodynamics (ClsElcDnm).
In Order to understand the physical cause for the origin of the relativistic characteristics of the DrEl and some special feature of its quantum behaviour we have to investigate the participate of its FnSpr ElmElcChrg in the IstThrDmnRltQnt Schrodinger fermion vortical harmonic oscillation motion (SchFrmVrtHrmOscMtn), describing the inner structure of the SchEl. Therefore we shall try in following to describe the IstThrDmnRltQnt SchFrmVrtHrmOscMtn of the DrEl’s FnSpr ElmElcChrg by means of the well-known mathematical apparatus of the RltQntMch and to give a green light of a new physical interpretation by virtue of the known language and conceptions of the NrlClsMch. We begin our new description by writing the well-known linear partial differential wave equation in partial derivative (LnrPrtDfrWvEqtPrtDrv) of Dirac, describing the RltQntMch behaviour of the DrEl. As it is well-known there exist different representations of the LnrPrtDfrWvEqtPrtDrv of Dirac within the RltQntMch. For example we begin with the presentation of its symmetrical representation :
$$i\mathrm{}\frac{\mathrm{\Psi }}{t}=H_d\mathrm{\Psi }=C(\alpha _j\widehat{p}_j)\mathrm{\Psi }+mC^2\beta \mathrm{\Psi };$$
(1)
where matrices $`\alpha _j`$ and $`\beta `$ have the well-known form :
$$\alpha _j=\left|\begin{array}{cc}\stackrel{~}{0}& \stackrel{~}{\sigma }_j\\ \stackrel{~}{\sigma }_j& \stackrel{~}{0}\end{array}\right|\mathrm{and}\beta =\left|\begin{array}{cc}\stackrel{~}{1}& \stackrel{~}{0}\\ \stackrel{~}{0}& \stackrel{~}{1}\end{array}\right|;$$
(2)
The total wave function (TtlWvFnc) $`\mathrm{\Psi }`$ of the DrEl within the NrlQntMch, satisfying the LnrPrtDfrWvEqtPrtDrv of Dirac (1), has four components $`(\psi _1,\psi _2,\psi _3,\psi _4)`$. Dirac had secured Lorentz’ invariation of his LnrPrtDfrWvEqtPrtDrv by dint of the introduction of four matrices $`(\alpha _1,\alpha _2,\alpha _3,\beta )`$.
There exist also the standard representation :
$$i\mathrm{}\frac{\stackrel{~}{\mathrm{\Psi }}}{t}=\stackrel{~}{H}_d\stackrel{~}{\mathrm{\Psi }}=C(\gamma _j\widehat{p}_j)\stackrel{~}{\mathrm{\Psi }}+mC^2\gamma _o\stackrel{~}{\mathrm{\Psi }};$$
(3)
where matrices $`\gamma _j`$ and $`\gamma _o`$ have the well-known form :
$$\gamma _j=\left|\begin{array}{cc}\stackrel{~}{\sigma }_j& \stackrel{~}{0}\\ \stackrel{~}{0}& \stackrel{~}{\sigma }_j\end{array}\right|\mathrm{and}\gamma _o=\left|\begin{array}{cc}\stackrel{~}{0}& \stackrel{~}{1}\\ \stackrel{~}{1}& \stackrel{~}{0}\end{array}\right|$$
(4)
In order to understand the physical meaning of the both TtlWvFnc $`\mathrm{\Psi }`$ and $`\stackrel{~}{\mathrm{\Psi }}`$ we must rewrite the both LnrPrtDfrWvEqnsPrtDvt of DrEl (1) and (3) in their component forms. In such a way we may write the following system of the motion equations :
$$i\mathrm{}\frac{\phi }{t}=C(\sigma _j\widehat{P}_j)\chi +mC^2\phi ;i\mathrm{}\frac{\chi }{t}=C(\sigma _j\widehat{P}_j)\phi mC^2\chi ;$$
(5)
$$i\mathrm{}\frac{\eta }{t}=C(\sigma _j\widehat{P}_j)\eta +mC^2\lambda ;i\mathrm{}\frac{\lambda }{t}=C(\sigma _j\widehat{P}_j)\phi +mC^2\eta ;$$
(6)
The eqs.(5) and (6) have been written by virtue of the following useful designations of the component sets of both the TtlWvFcs $`\mathrm{\Psi }`$ and $`\stackrel{~}{\mathrm{\Psi }}`$ :
$$\phi =\left|\begin{array}{c}\psi _1\\ \psi _2\end{array}\right|;\mathrm{and}\chi =\left|\begin{array}{c}\psi _3\\ \psi _4\end{array}\right|;\mathrm{or}\eta =\left|\begin{array}{c}\stackrel{~}{\psi }_1\\ \stackrel{~}{\psi }_2\end{array}\right|;\mathrm{and}\lambda =\left|\begin{array}{c}\stackrel{~}{\psi }_3\\ \stackrel{~}{\psi }_4\end{array}\right|;$$
(7)
As it is following from (5) the total energy of the DrEl at its forward motion $`i\mathrm{}\frac{\phi }{t}`$ is equal of the sum of its kinetic energy $`C(\stackrel{~}{\sigma }_j\stackrel{~}{P}_j)`$ in the state $`,\chi `$ and potential energy $`m.C^2`$ of its ElcMgnSlfAct of its FnSpr ElmElcFld with a corresponding resultant own QntElcMgnFld, created by the StchVrtPhtns, radiated by its FnSpr ElmElcChrg in the state $`\phi `$. Just in such a way about the total energy of the DrEl $`i\mathrm{}\frac{\chi }{t}`$ is equal of the sum of its kinetic energy $`C(\stackrel{~}{\sigma }_j\stackrel{~}{P}_j)`$ in the state $`\phi `$ and potential energy $`m.C^2`$ of its ElcMgnSlfAct of its FnSpr ElmElcFld with a corresponding resultant own QntElcMgnFld, created by the StchVrtPhtns, radiated by its FnSpr ElmElcChrg in the state $`\chi `$. That is because we may flatly assert that the components $`\phi `$ of its TtlWvFnc $`\mathrm{\Psi }`$ in the symmetrical presentation describe the forward motion of the DrEl and the components $`\chi `$ of the same TtlWvFnc describe the backward motion of the DrEl. Besides that if the both odd components: $`\mathrm{\Psi }_1`$ and $`\mathrm{\Psi }_3`$ describe the DrEl’s spinning in a left, then both even components : $`\mathrm{\Psi }_2`$ and $`\mathrm{\Psi }_4`$ describe the DrEl’s spinning in a right. Therefore these two group WvFnc have opportunity sign before the angle $`\phi `$ are transforming separately.
The existence relations between both TtlWvFnc $`\mathrm{\Psi }`$ and $`\stackrel{~}{\mathrm{\Psi }}`$ give us a clear physical explanation and correct mathematical description of all components of the DrEl’s TtlWvFnc in the both representation. From the LnrPrtDfrWvEqtPrtDrv of Dirac we may see that matrixes work within its as switches, making possible the correct mathematical description of the IstThrDmnRltQnt SchFrmVrtHrmOscMtn of the DrEl’s FnSpr ElmElcChrg by dint of the four components of its TtlWvFnc.
In a due course it is easily to show further that the proper mechanical moment (PrpMchMmn)(spin) of the DrEl can really be created as a result of the participate of its FnSpr ElmElcChrg in the IstThrDmnRltQnt SchFrmVrtHrmOscMtn. In order to obtain this in a naturally way we have to rewrite the four one-component LnrPrtDfrWvEqtPrtDrv of Dirac (5) in more obvious form :
$`i\mathrm{}{\displaystyle \frac{\psi _1}{t}}+i\mathrm{}C\times \left\{{\displaystyle \frac{\psi _4}{x}}i{\displaystyle \frac{\psi _4}{y}}+{\displaystyle \frac{\psi _3}{z}}\right\}=mC^2\psi _1;`$
$`i\mathrm{}{\displaystyle \frac{\psi _2}{t}}+i\mathrm{}C\times \left\{{\displaystyle \frac{\psi _3}{x}}+i{\displaystyle \frac{\psi _3}{y}}{\displaystyle \frac{\psi _4}{z}}\right\}=mC^2\psi _2;`$
$`i\mathrm{}{\displaystyle \frac{\psi _3}{t}}+i\mathrm{}C\times \left\{{\displaystyle \frac{\psi _2}{x}}i{\displaystyle \frac{\psi _2}{y}}+{\displaystyle \frac{\psi _1}{z}}\right\}=mC^2\psi _3;`$
$`i\mathrm{}{\displaystyle \frac{\psi _4}{t}}+i\mathrm{}C\times \left\{{\displaystyle \frac{\psi _1}{x}}+i{\displaystyle \frac{\psi _1}{y}}{\displaystyle \frac{\psi _2}{z}}\right\}=mC^2\psi _4;`$ (8)
Then the system of four one-component LnrPrtDfrWvEqtPrtDrv of Dirac (5) may be rewritten by means of the substitution : $`x=\rho \mathrm{cos}\varphi `$ and $`y=\rho \mathrm{sin}\varphi `$; from Decart’s coordinates in more obvious form of cylindrical coordinates :
$`i\mathrm{}{\displaystyle \frac{\psi _1}{t}}+i\mathrm{}C\left[\mathrm{exp}i\varphi \left\{{\displaystyle \frac{\psi _4}{\rho }}{\displaystyle \frac{i}{\rho }}{\displaystyle \frac{\psi _4}{\varphi }}\right\}+{\displaystyle \frac{\psi _3}{z}}\right]=mC^2\psi _1;`$
$`i\mathrm{}{\displaystyle \frac{\psi _2}{t}}+i\mathrm{}C\left[\mathrm{exp}+i\varphi \left\{{\displaystyle \frac{\psi _3}{\rho }}+{\displaystyle \frac{i}{\rho }}{\displaystyle \frac{\psi _3}{\varphi }}\right\}+{\displaystyle \frac{\psi _4}{z}}\right]=mC^2\psi _2;`$
$`i\mathrm{}{\displaystyle \frac{\psi _3}{t}}+i\mathrm{}C\left[\mathrm{exp}i\varphi \left\{{\displaystyle \frac{\psi _2}{\rho }}{\displaystyle \frac{i}{\rho }}{\displaystyle \frac{\psi _2}{\varphi }}\right\}+{\displaystyle \frac{\psi _1}{z}}\right]=mC^2\psi _3;`$
$`i\mathrm{}{\displaystyle \frac{\psi _4}{t}}+i\mathrm{}C\left[\mathrm{exp}+i\varphi \left\{{\displaystyle \frac{\psi _1}{\rho }}+{\displaystyle \frac{i}{\rho }}{\displaystyle \frac{\psi _1}{\varphi }}\right\}{\displaystyle \frac{\psi _2}{z}}\right]=mC^2\psi _4;`$ (9)
As it is easily to seen if the first pair of TtlWvFnc’s components $`\psi _1`$ and $`\psi _3`$ have equal phase factors $`\varphi _1`$, than the second the pair of TtlWvFnc’s components $`\psi _2`$ and $`\psi _3`$ have also equal phase factors $`\varphi _1`$. As it follows from eq.(9) the difference between two phase factors is equal of $`\varphi `$. Therefore we may suppose by means of a symmetrical consideration that four components (OrbWvFnc) of the DrEl’s TtlWvFnc $`\mathrm{\Psi }`$ have the following presentations :
$`\psi _1(\rho ,\varphi ,z)=\overline{\psi }_1(\rho ,\varphi ,z)\mathrm{exp}i(\varphi /2);\psi _3(\rho ,\varphi ,z)=\overline{\psi }_3(\rho ,\varphi ,z)\mathrm{exp}i(\varphi /2);`$ (10)
$`\psi _2(\rho ,\varphi ,z)=\overline{\psi }_2(\rho ,\varphi ,z)\mathrm{exp}+i(\varphi /2);\psi _4(\rho ,\varphi ,z)=\overline{\psi }_4(\rho ,\varphi ,z)\mathrm{exp}+i(\varphi /2);`$ (11)
If we take into account that the participate of the well spread (WllSpr) ElmElcChrg of the SchEl in the IstThrDmnRltQnt FrthStchBsnCrcHrmOscMtn, securing its quantum behavior secures an additional dispersion $`(\delta \varphi /2)`$, then the TtlWvFnc $`\mathrm{\Psi }`$ of the DrEl can be rewritten in two following representations :
$`\mathrm{\Psi }_{l+1/2}(\rho ,\varphi ,z)={\displaystyle \frac{\varphi _o}{2}}\left|\begin{array}{c}\psi _{1l}(\rho ,z)\mathrm{exp}il\varphi \\ \psi _{2l}(\rho ,z)\mathrm{exp}i(l+1)\varphi \\ \psi _{3l}(\rho ,z)\mathrm{exp}il\varphi \\ \psi _{4l}(\rho ,z)\mathrm{exp}i(l+1)\varphi \end{array}\right|`$ (16)
and
$`\mathrm{\Psi }_{l1/2}(\rho ,\varphi ,z)={\displaystyle \frac{\varphi _o}{2}}\left|\begin{array}{c}\psi _{1l}(\rho ,z)\mathrm{exp}i(l1)\varphi \\ \psi _{2l}(\rho ,z)\mathrm{exp}il\varphi \\ \psi _{3l}(\rho ,z)\mathrm{exp}i(l1)\varphi \\ \psi _{4l}(\rho ,z)\mathrm{exp}il\varphi \end{array}\right|`$ (21)
In the meanwhile it is easily to verify by virtue of the operator
$$\widehat{J}_z=\left\{i\mathrm{}\frac{}{\varphi }+\frac{\mathrm{}}{2}\sigma _z\right\}$$
(22)
that if the TtiWvFnc (16) describes the behaviour of the free DrEl, having the TtiMchMmn’s value $`J_z=\mathrm{}(l+1/2)`$, than the TtiWvFnc (21) describes the behaviour of the free DrEl, having the TtiMchMmn’s value $`J_z=\mathrm{}(l1/2)`$. Indeed, if
$`\widehat{J}_z\mathrm{\Psi }_{l+1/2}(\rho ,\varphi ,z)={\displaystyle \frac{\varphi _o}{2}}\left|\begin{array}{c}\mathrm{}(l+1/2)\psi _{1l}(\rho ,z)\mathrm{exp}il\varphi \\ \mathrm{}(l+11/2)\psi _{2l}(\rho ,z)\mathrm{exp}i(l+1)\varphi \\ \mathrm{}(l+1/2)\psi _{3l}(\rho ,z)\mathrm{exp}il\varphi \\ \mathrm{}(l+11/2)\psi _{4l}(\rho ,z)\mathrm{exp}i(l+1)\varphi \end{array}\right|=\mathrm{}(l+1/2)\mathrm{\Psi }_{l+1/2}(\rho ,\varphi ,z)`$ (27)
and
$`\widehat{J}_z\mathrm{\Psi }_{l1/2}(\rho ,\varphi ,z)={\displaystyle \frac{\varphi _o}{2}}\left|\begin{array}{c}\mathrm{}(l1+1/2)\psi _{1l}(\rho ,z)\mathrm{exp}i(l1)\varphi \\ \mathrm{}(l1/2)\psi _{2l}(\rho ,z)\mathrm{exp}il\varphi \\ \mathrm{}(l1+1/2)\psi _{3l}(\rho ,z)\mathrm{exp}i(l1)\varphi \\ \mathrm{}(l1/2)\psi _{4l}(\rho ,z)\mathrm{exp}il\varphi \end{array}\right|=\mathrm{}(l1/2)\mathrm{\Psi }_{l1/2}(\rho ,\varphi ,z)`$ (32)
The presented upper investigation shows us that the FnSpr ElmElcChrg of the DrEl really participates in the IstThrDmnRltQnt SchFrmVrtHrmOscMtn and the WllSpr ElmElcChrg of the SchEl really participates in the IstThrDmnRltQnt FrthStchBsnCrcHrmOscMtn. In what following we wish to show that the participating the FnSpr ElmElcChrg in its self-consistent IstThrDmnRltQnt SchFrmVrtHrmOscMtn (Zitterbevegung) is very effective at creation of its own resultant self-consistent quantized electromagnetic field (QntElcMgnFld) by means of stochastic emission of the VrtPhtns.
In the first Breit and afterwards Fock had observed that the instantaneous velocity operator $`\frac{\widehat{r}_j}{t}`$ assume, that the operator of the instantaneous velocity of a free QntMcrPrt have very paradoxical form in the RltQntMch. Indeed, it is well-known from the RltQntMch that the analytical operator form of the instantaneous velocity value of the free DrEl may be obtained by virtue of the Heisenberg commutation relations between the operators of its radius-vector $`\widehat{r}_j`$ and Dirac’s hamiltonian $`H_d`$. In such the way we can obtain:
$$\widehat{V}_j=\frac{d\widehat{r}_j}{dt}=\frac{i}{\mathrm{}}(\overline{r}_j\overline{H}_d\overline{H}_d\overline{r}_j)=C\alpha _j;$$
(33)
Seventy years ago Schrodinger had investigated the physical meaning of the operators in the RltQntMch, describing the relativistic quantized behaviour of the DrEl in the old Dirac picture, making use of its linear Hamiltonian, LnrPrtDfrWvEqtPrtDrv of Dr and four component TtlWvFnc $`\mathrm{\Psi }(r,t)`$. First of all he had obtained the motion equation for Dirac’s matrices :
$$i\mathrm{}\frac{\alpha _j}{t}=(\alpha _jH_dH_d\alpha _j)=\mathrm{\hspace{0.17em}2}(\alpha _jH_dC\widehat{p}_j)=\mathrm{\hspace{0.17em}2}(C\widehat{p}_j\alpha _jH_d);$$
(34)
After replacing in the equation (34) the matrices $`\alpha _j`$ by the matrices $`\eta _j`$ according to the following equation :
$$\eta _j=\left(\alpha _j\frac{C\widehat{p}_j}{H_d}I_o\right);$$
(35)
Schrodinger had obtained the oscillation equation for $`\eta `$ matrices. In such a way he had obtained the following solution of eq.(33) for $`r(t)_j`$ :
$$\widehat{r}_j=\widehat{a}_j+\frac{tC^2\widehat{p}_j}{H_d}I_o+\frac{iC\mathrm{}}{H_d}\left(\alpha _j\frac{C\widehat{p}_j}{H_d}I_o\right)_{t=o}\mathrm{exp}\{\frac{2itH_d}{\mathrm{}}\};$$
(36)
From eq.(36) we can see that the second term describes the classical motion of the free ClsMcrPrt and one increases in a linear way with the current velocity when the time is increasing. The last term $`\eta `$ in the eq.(36) describes the self-consistent inner motion of the FnSpr ElmElcChrg of the DrEl, which had called Zitterbewegung by Schrodinger. As would be understand the participation of the SchEl’s WllSpr ElmElcChrg in the IstThrDmnNrlQnt FrthStchBsnCrcHrmOscsMtm is not described by its OrbWvFnc and therefore there are no own part in the coordinate operator (36). The participation of the DrEl’s FnSpr ElmElcChrg in the IstThrDmnRltQnt SchFrmVrtHrmOscMtn, which is well described by the Dirac’s matrices $`\alpha _j`$ and the four components of the DrEl’s total wave function (TtlWvFnc), directs us to construct a new matrix’ RltQntElcDnm in an accordance with the Maxwell’s nonrelativistic ClsElcDnm, where the classical motion of some NtnClsPrt is described with the smooth thin line. Moreover, as it will be obvious in the following investigation there exists a possibility to understand the physical reasons of the rest self-energy origin and the physical interpretation of the LnrPrtDfrWvEqt of Dirac.
Indeed, we shall demonstrate in the what following that because of the participation of the DrEl’s FnSpr ElmElcChrg in its IstThrDmnRltQnt SchFrmVrtHrmOsc motion, there exists a possibility to calculate also the instantaneous RslSlfCns values of all the components of the ElcInt $`E_j`$ and of the MgnInt $`H_j`$ of the own resultant QntElcMgnFld in the point of its instantaneous position by means of the Dirac’s matrices, describing this self-consistent motion. As it was shown above the RslQwn QntElcMgnFld is begotten by Hgh-Enr StchVrtPhtns, emitted by its FnSpr ElmElcChrg in different time moments from corresponding different space positions, being occupied by its FnSpr ElmElcChrg during the last half-period.
Indeed, for long time (about one centaur) ago, thence R.Schwarzschild had written the electro-kinetic term $`\rho (\phi v.A)`$ into the Lagranjian density, it is well-known from the Maxwell ClsElcDnm that when the LrEl is found in the external ClsElcMgnFld its canonical impulse $`p_j`$ amounts in two parts :
$$P_j=p_j^k+p_j^p=mv_j\frac{e}{C}A_j;W=E\frac{e}{C}\phi ;$$
(37)
The first part $`p_j^k=mv_j`$ is the well-known kinematic momentum and the second part $`p_j^p=\frac{e}{C}A_j`$ is the potential momentum, which the DrEl with his FnSpr ElmElcChrg have acquired when it was brought in some external ClsElcMgnFld. Therefore the uncommon behaviour of the SchEl in the external QntElcMgnFld may be described as the behavior of the free one, only replacing the canonical impulse operator $`\widehat{p}_j=i\mathrm{}_j`$ with its generalized impulse $`P_j`$ and generalized energy $`W`$, described in the following form (37):
$$\widehat{P}_j=\widehat{p}_j\frac{e}{C}\widehat{A}_j;\widehat{W}=\widehat{E}\frac{e}{C}\widehat{\phi };$$
(38)
with canonical impulse operator $`\widehat{P}_j=i\mathrm{}_j`$ and canonical energy operator $`\widehat{W}=i\mathrm{}\frac{}{t}`$. For the sake of the intrinsic symmetry it is convenient to make use of the generalized impulse of its zero component $`P_o`$ instead of $`\frac{\widehat{W}}{C}`$ and of the mechanical impulse of its zero component $`p_o`$ instead of $`\frac{\widehat{E}}{C}`$. However it is easy to verify that if all the components of the canonical impulses commute between them-selves :
$`\widehat{P}_j\widehat{P}_l\widehat{P}_l\widehat{P}_j=\delta _{jl};and\widehat{P}_j\widehat{P}_o\widehat{P}_o\widehat{P}_j=\delta _{jo};`$ (39)
then all the components of the kinetic impulse don’t commute between oneself. It is turned out that the commutations between the mechanical (kinematic) impulse generate the values of the ElcInt $`E_j`$ and the MgnInt $`H_j`$ of the external ClsElcMgnFld, as it is well-known from the NrlQntMch. Indeed,if we rewrite the eqs.$`(\text{38})`$ in the following form :
$$i\mathrm{}_j=\widehat{p}_j\frac{e}{C}\widehat{A}_j=m\widehat{v}_j\frac{e}{C}\widehat{A}_j;$$
(40)
and
$$i\frac{\mathrm{}}{C}\frac{}{t}=\widehat{p}_o\frac{e}{C}\widehat{A}_o=\frac{\widehat{E}}{C}\frac{e}{C}\widehat{\phi };$$
(41)
then the four commutations between the four mechanical (kinematic) impulse components generate the six values of the ElcInt $`E_j`$ and of the MgnInt $`H_j`$ of the external ClsElcMgnFld, as it is well-known from the NrlQntMch :
$$\widehat{p}_j\widehat{p}_l\widehat{p}_l\widehat{p}_j=m\widehat{v}_jm\widehat{v}_lm\widehat{v}_lm\widehat{v}_j=i\mathrm{}\frac{e}{C}\epsilon _{jlk}H_k;$$
(42)
and
$$\widehat{p}_j\widehat{p}_o,\widehat{p}_o\widehat{p}_j=m\widehat{v}_jmhatv_om\widehat{v}_om\widehat{v}_j=i\mathrm{}\frac{e}{C}E_j;$$
(43)
The six components values of the ElcInt $`E_j`$ and the MgnInt $`H_j`$ of the external QntElcMgnFld are determined for the space position $`r_j`$ of the WllSpr ElmElcChrg of the SchEl (from the NrlClsMch’s point of view) at the time moment t.Indeed, there aren’t any mistake in our physical interpretation. Indeed, we have to remember that the IstThrDmnNrlQnt FrthStchBsnCrcHrmOscs’ motion is a result of the permanent interaction of the WllSpr ElmElcChrg of the SchEl with the ElcInt $`E_j`$ of the StchVrtPhtns, generated in the neighbour area of its localization by the FlcVcm. Hence as the IstThrDmnNrlQnt FrthStchBsnCrcHrmOscs’motion of the QntMcrPrt has no a thin and smooth classical trajectory, which may be determined as a solution of the Newton equation, then we are forced to take advantage of results of the Heisenberg commutations (42) and (43). Moreover we have no right to take into account the screening influence of the FlcVcm over the ElcInt at the description of the ElcMgmSlfInt of the FnSpr ElmElcChrg and MgnDplMmn of the DrEl with the RslSlfCns values of the ElcInt and the MgnInt of own resultant QntElcMgnFld as we already use it at taking into account its influence at the description of the SchEl’s participation in the IstThrDmnNrlQnt FrthStchBsnCrcHrmOsc’s motion.
Consequently, the vector-potential $`\stackrel{ˇ}{A}_j`$ and potential $`\stackrel{ˇ}{V}_j`$, which takes part within the LnrPrtDfrWvEqt of Dirac don’t contain any contribution of the existent StchVrtPhtns within the FlcVcm. However on other hand it is known from the RltQntMch too that the operator of the instantaneous velocity $`V_j`$ of the free FnSpr DrEl has the exceptional form $`C\alpha _j`$ as it is seen from eq.(33). Therefore if we want to describe correctly the IstThrDmnRltQnt SchFrmVrtHrmOsc’s motion of the free DrEl, who is moving by its instantaneous velocity $`\widehat{V}_j=C\alpha _j`$ within own resulting QntElcMgnFld by virtue of the mathematical apparatus of the RltQntMch, we must take use of the well-known Heisenberg commutation relations (42) and (43). Therefore we may suppose that the DrEl’s generalized moments in the RltQntMch could be determined by Dirac’s matrices as they ware done the velocity was exchange by $`C\alpha _j`$ in the same eqs.(36). Then we have obtain the following presentation of the DrEl’s mechanical momentum components in the motionless coordinate system of reference relatively to the centre of the space distribution of the SchEl’s WllSpr ElmElcChrg :
$$\stackrel{ˇ}{p}_o=mC\widehat{I};\mathrm{and}\stackrel{ˇ}{p}_j=mC\alpha _j;$$
(44)
where $`\stackrel{ˇ}{I}`$ denotes the unitar matrix. Hence the HsnCmtRlt between the kinetic momentum components (44) of the DrEl in the RltQntMch, which are analogous of the HsnCmtRlt for the SchEl to eqs.(42) and (43). have to determine the RslSlfCns values of the ElcInt $`\stackrel{ˇ}{E}_j`$ and the MgnInt $`\stackrel{ˇ}{H}_j`$ of own QntElcMgnFld. Therefore they may be written in the following form :
$$\stackrel{ˇ}{p}_j\stackrel{ˇ}{p}_l\stackrel{ˇ}{p}_l\stackrel{ˇ}{p}_j=m^2C^2(\alpha _j\alpha _l\alpha _l\alpha _j)=\mathrm{\hspace{0.17em}2}im^2C^2\epsilon _{jlk}\sigma _k;$$
(45)
and
$$\stackrel{ˇ}{p}_j\stackrel{ˇ}{p}_o\stackrel{ˇ}{p}_o\stackrel{ˇ}{p}_j=m^2C^2(\alpha _j\widehat{I}\widehat{I}\alpha _j)=\mathrm{\hspace{0.17em}0};$$
(46)
if the commutations (42) and (43) between the mechanical (kinematic) impulse $`\widehat{p}_j,\widehat{p}_l`$ and $`\widehat{p}_o`$ generate the AvrVls $`\widehat{E}_j`$ and $`\widehat{H}_j`$ of the ElcInt and MgnInt of the external QntElcMgnFld, then the commutations (45) and (46) between the mechanical (kinematic) impulse $`\stackrel{ˇ}{p}_j,\stackrel{ˇ}{p}_l`$ and $`\stackrel{ˇ}{p}_o`$ generate the PslSlfCnsVls $`\stackrel{ˇ}{E}_j`$ and $`\stackrel{ˇ}{H}_j`$ of the ElcInt and the MgnInt of the own resultant QntElcMgnFld, which may be obtained by means of the comparison of the corresponding right parts of the (42), (43) and (45), (46) and one may be described by the following formulas :
$$\stackrel{ˇ}{E}_j=\mathrm{\hspace{0.17em}0};\mathrm{and}\stackrel{ˇ}{H}_j=\mathrm{\hspace{0.17em}2}\frac{m^2C^3}{e\mathrm{}}\sigma _j;$$
(47)
In such an easy way only owing to a supposition of the self-consistency of IstThrDmnRltQnt SchFrmVrtHrmOsc’s motion of the DrEl’s FnSpr ElmElcChrg at about the light velocity C, who minimizes in a self-consistent way the self-energy at a rest of its electromagnetic self-acting (ElcMgnSlfAct) between its continuously moving FnSpr ElmElcChrg and the corresponding electric current with their corresponding potential and vector-potential. The RslSlfCnsVls of own QntElcMgnFld of the DrEl is created by its FnSpr ElmElcChrg, emitting incessantly high energy StchVrtPhtns at different moments of the latest time from its various corresponding positions in a space and absorbed in the form of the ElcMgnSlfAct by the DrEl’s FnSpr ElmElcChrg in its instantaneous positions. All the RslSlfCnsVls of the ElcInt components $`\stackrel{ˇ}{E}_j`$ of own QntElcMgnFld may be precisely compensated in a result of the participation of the DrEl’s FnSpr ElmEl in same IstThrDmnRltQnt SchFrmVrtHrmOscsMtn. Only in a result of such the IstThrDmnRltQnt SchFrmVrtHrmOsc’s motion all the RslSlfCns values of the MgnInt components $`\stackrel{ˇ}{H}_j`$ of own QntElcMgnFld may obtain double values in a comparison with the corresponding averaged values of the MgnInt components $`\widehat{H}_j`$ of the ClsElcMgnFld of some NtnClsPrt, charged by the FnSpr ElmElcChrg, fulfilled the IstThrDmnRltQnt BsnStchCrcHrmOsc’s motion in a conformity with the laws of the Einstein relativistic classical mechanics (RltClsMch).
The Schrodinger’s zitterbewegung is a self-consistent strong correlated vortical motion, which minimizes the self-energy at a rest of the QntMcrPrt and one secures the continuous stability of its Shrodinger’s wave package (SchWvPct) in time within the space. Therefore this inner self-consistent quantized vortical motion of the QntMcrPrt’s FnSpr ElmElcChrg corresponds to its inner harmonical motion, introduced firstly by Louis de Broglie.
There exists some possibilities enough not only for rough calculations of the averaged undivergent potential and vector-potential of the DrEl’s FnSpr ElmElcChrg and real values of its particular MgnDplMm $`\mu _j`$ and MchMmn $`s_j`$, which are results of the participation of the DrEl’s FnSpr ElmElcChrg in its IstThrDmnRltQnt SchFrmVrtHrmOsc’s motion. This natural conclusion explains the physical cause of the doubled gyromagnetical ratio of the DrEl’s inner MgnDplMmn $`\mu _j=\frac{e\mathrm{}}{2mC}\sigma _j`$ to its inner MchMmn (spin) $`S_j=\frac{\mathrm{}}{2}\sigma _j`$ and abolishes the necessity of useless renormalization of its ElcChrg and mass because of the absence of any physical substantiations, as will be seen in further researches. Hence the magnetic productivity of the DrEl’s FnSpr ElmElcChrg as a result of its participation in the IstThrDmnRltQnt SchFrmVrtHrmOsc’s motion exceeds in twice the magnetic productivity of the SchEl’s WllSpr ElmElcChrg as a result of its participation of IstThrDmnNrtQnt FrthStchBsnHrmOsc’s motion with same parameters.
In the proposed PhsMdl of the DrEl its rest energy $`m.C^2`$ may be considered as a natural consequence to the ElcMgnSlfAct between its RslSlfCns values of the MgnInt $`\stackrel{ˇ}{H}_j=\mathrm{\hspace{0.17em}2}(\frac{m^2C^3}{e\mathrm{}})\sigma _j`$ of the own QntElcMgnFld in the point of its instantaneous position and the RslSlfCns values of its MgnDplMmn $`\mu _j=(\frac{e\mathrm{}}{2mC})\sigma _j`$ at same point :
$$E_o=\mu _j\stackrel{ˇ}{H}_j=\frac{e\mathrm{}}{2mC}\sigma _j\frac{2m^2C^3}{e\mathrm{}}\sigma _j=mC^2\sigma _j\sigma _j=mC^2;$$
(48)
Indeed, it is easy to verify that :
$$\sigma _j\sigma _j=\mathrm{sin}^2\theta cos^2\varphi +\mathrm{sin}^2\theta sin^2\varphi +cos^2\theta =\mathrm{\hspace{0.17em}1};$$
(49)
As a result of the above investigation we can affirm that the participation of the SchEl’s WllSpr ElmElcChrg in IstThrDmnNrlQnt FrthStchCrcBsnHrmOscMtn cause existence of its anomalous MgnDplMmn $`\delta \mu _o`$. Therefore if we wish to obtain a ratio of the anomalous MgnDplMmn $`\delta \mu _o`$ of the SchEl’s WllSpr ElmElcChrg as a result of its participate in the IstThrDmnNrlQnt FrthStchCrcBsnHrmOscMtn to its own MgnDplMmn $`\mu _o`$ as a result of the participation of DrEl’s FnSpr ElmElcChrg in the IstThrDmnNrlQnt SchFrmHrmOscMtn, then we must know that it is equal to half ratio of their kinetic energies. Consequently as a result of the natural relations we can obtain :
$$\frac{\delta \mu _o}{\mu _o}=\frac{1}{2\pi }\frac{e^2}{C\mathrm{}}\frac{mC^2}{mC^2};and\delta \mu _o=\frac{\mu _o}{2\pi }\frac{e^2}{C\mathrm{}};$$
(50)
In a consequence of the above used approach we may propose that the RslSlfCns values of the four components of the vector-potential $`\stackrel{ˇ}{A}_j`$ and potential $`\stackrel{ˇ}{\varphi }`$ of the DrEl’s FnSpr ElmElcChrg, participating in the IstThrDmnNrlQnt SchFrmVrtHrmOsc’s motion, one have the following analytic form :
$$\stackrel{ˇ}{A}_j=\frac{m.C^2}{e}\mathrm{\Psi }^{}|\beta \alpha _j|\mathrm{\Psi };\mathrm{and}\stackrel{ˇ}{\varphi }=\frac{m.C^2}{e}\mathrm{\Psi }^{}|\beta |\mathrm{\Psi };$$
(51)
There exists a physical cause and a mathematical possibility for obtaining the RslSlfCns values $`(\text{47})`$ of the ElcInt $`\stackrel{ˇ}{E}_j`$ and the MgnInt $`\stackrel{ˇ}{H}_j`$ from the four components in eqs.(51) of their vector-potential $`\stackrel{ˇ}{A}_j`$ and potential $`\stackrel{ˇ}{\varphi }`$ by means of the corresponding relations between analogous ones of that, expressing the Maxwell’s laws within the Maxwell’s ClsElcDnm. For that purpose we have to exchange the coordinate operators $`_j`$ and time operator $`\frac{1}{C}\frac{}{Ct}`$ in the Maxwell’s ClsElcDnm’s scheme, where the ElcChrg McrPrt’s motion is described by its coordinate $`r_j`$, by the matrix operators $`i\frac{mC}{\mathrm{}}\alpha _j`$ and $`\frac{mC}{\mathrm{}}\widehat{I}`$, when the FnSpr ElmElcChrg’s motion is described by the matrices, in accordance with the relations (42), (43), (45), (46), and (47). Indeed, then instead of the Maxwell’s relations :
$$H_j=\left[\times A\right]_j=\epsilon _{jkl}_kA_l;$$
(52)
and
$$E_j=\frac{1}{C}\frac{A_j}{t}_j\varphi $$
(53)
we must write the following unusual equations :
$$\stackrel{ˇ}{H}_j=i\epsilon _{jkl}\frac{mC}{\mathrm{}}\alpha _j\frac{mC^2}{e}\alpha _l=\mathrm{\hspace{0.17em}2}\frac{m^2C^3}{e\mathrm{}}\sigma _j;$$
(54)
and
$$\stackrel{ˇ}{E}_j=i\frac{mC}{\mathrm{}}\widehat{I}\frac{mC^2}{e}\alpha _ji\frac{mC}{\mathrm{}}\alpha _j\frac{mC^2}{e}\widehat{I}=\mathrm{\hspace{0.17em}0};$$
(55)
There is a necessity to note here that the presence of the matrices in the (47) is connected with the physical reason that the moment RslSlfCns values of the potential and vector-potential in the time $`t_1`$ are results of the interference of the QntElcMgnFlds of all Hgh-Enr StchVftPhtns, emitted in the time interval $`(t_1\frac{T}{2})tt_1`$. It is easy to be shown in Veile’s symmetrical representation, where the matrices $`\gamma _j`$ and $`\gamma _o`$ play the parts of the matrices $`\alpha _j`$ and $`\beta `$, participating in the Pauli-Dirac’s representation.
As a consequence of our felicitious supposition, using by me at building of our PhsMdl, we may write the LnrDfrWvEqnPrtDrv of Dirac by means of the expression (51) for the RslSlfCns values of the own potential $`\stackrel{ˇ}{\varphi }`$ and vector-potential $`\stackrel{ˇ}{A}_j`$ of the DrEl’s FnSpr ElmElcChrg and its corresponding current in the following form :
$`\mathrm{\Psi }^+,|H_d|\mathrm{\Psi }=e\mathrm{\Psi }^+|\mathrm{\Psi }{\displaystyle \frac{m.C^2}{e}}\mathrm{\Psi }^+|\beta |\mathrm{\Psi }+`$
$`eC\mathrm{\Psi }^+|\alpha _j|\mathrm{\Psi }{\displaystyle \frac{mC}{e}}\mathrm{\Psi }^+|\beta \alpha _j|\mathrm{\Psi }+C\mathrm{\Psi }^+|\alpha _jP_j|\mathrm{\Psi };`$ (56)
Then if we take into account the existence of the following equations :
$$\mathrm{\Psi }^+|\mathrm{\Psi }=\mathrm{\hspace{0.17em}1};\mathrm{and}\mathrm{\Psi }^+|\alpha _j|\mathrm{\Psi }\mathrm{\Psi }^+|\beta \alpha _j|\mathrm{\Psi }=\mathrm{\hspace{0.17em}0};$$
(57)
then the LnrPrtDfrWvEqn of Dr. The (2) may be rewritten into the following well-known form :
$$\mathrm{\Psi }^+|H_d|\mathrm{\Psi }=C\mathrm{\Psi }^+|\alpha _jP_j|\mathrm{\Psi }+mC^2\mathrm{\Psi }^+|\beta |\mathrm{\Psi }$$
(58)
Our investigation shows the role of each component of the TtlWvFnc $`\mathrm{\Psi }`$ of the DrEl at the description of its behaviour and the role of the matrixes as useful switches, making possible the description of the DrEl’s motion by means of its four-component TtlWvFnc. Because of such our interpretation of the space-time dependent part of the SchEl’s one-component OrbWvFnc $`\mathrm{\Psi }`$ , which gives the description of its orbital motion only, one could be compared with black and white TV, while the space-time dependent part of the DrEl’s four-components TtlWvFnc $`\mathrm{\Psi }`$, which together with matrices give the description of its full motion,one should be compared with color TV. In other words, while one-component OrbWvFnc $`\psi `$ give us the colourless description only of the SchEl’s motion, while the four-components TtlWvFnc $`\mathrm{\Psi }`$ together with matrices give diversiform of the coloured description of the DrEl’s motion.
The existence of the well-known relativistic relation $`E^2=P^2C^2+m^2C^4`$ between the energy, momentum and mass of a real elementary particle (ElmPrt) also may be obtained by the Maxwell equations of the ClsElcDnm, taking into consideration the MgnInt between the momentum of the charged ClsMcrPrt and the vector-potential of its ClsElcMgnFld, participating in the IstThrDmnNrlQnt StchBsnVrtHrmOscMtn
In the long run I cherish hope that our consideration of the massive leptons’ behaviour from a new point of view of my felicitous PhsMdl, physically obvious expatiating the physical cause for the origin of all their properties, will be of great interest for all scientists. Theoretical physicists will find the picturesque explanation of all their properties and Philosophers will find the total employment of the dialectical materialism for the investigations of the micro particle’s world. My very quality and profound interpretation of the SchEl’s behaviour in the NrlQntMch and of the DrEl’s behaviour in the RltQntMch , as well as the excellent quantity coincidence of the deduced values of all DrEl’s parameters with their corresponding experimentally determined values gives as the hope for correctness of our beautiful, simple and preposterous PhsMdl and fine, extraordinary ideas, which have been inserted at its construction. |
warning/0002/cond-mat0002035.html | ar5iv | text | # Effect of Magnetic Impurities on Suppression of the Transition Temperature in Disordered Superconductors
## I Introduction
Many experiments performed on homogeneous disordered thin film superconductors have shown that superconductivity is suppressed by increasing disorder, as measured by the normal state resistance per square, $`R_{\text{ }\text{ }\text{ }\text{ }\text{ }}`$. The majority of the data is of the transition temperature as a function of the resistance per square, $`T_c(R_{\text{ }\text{ }\text{ }\text{ }\text{ }})`$, although there is some data on the upper critical field, $`H_{c2}(T,R_{\text{ }\text{ }\text{ }\text{ }\text{ }})`$, and the order parameter $`\mathrm{\Delta }(R_{\text{ }\text{ }\text{ }\text{ }\text{ }})`$. The main data to be explained thus consists of $`T_c(R_{\text{ }\text{ }\text{ }\text{ }\text{ }})`$ curves for different materials. To see how disorder might affect $`T_c`$, consider the mean field equation,
$$T_{c0}=1.13\omega _D\mathrm{exp}\left[\frac{1}{N(0)(\lambda \mu ^{})}\right],$$
(1)
where $`\lambda `$ is the attractive BCS interaction mediated by phonons of energy less than the Debye frequency, $`\omega _D`$, $`\mu ^{}`$ is the Coulomb pseudopotential, the effective strength of the Coulomb repulsion, and $`N(0)`$ is the single particle density of states at the Fermi surface. Obviously disorder could affect $`T_c`$ by changing $`\lambda `$, $`\mu ^{}`$, and $`N(0)`$. In particular the diffusive motion of electrons caused by the disorder is known to lead to an increased effective strength of the Coulomb interaction, as the screening is less efficient than with ballistic electrons, and this leads to an increase in $`\mu ^{}`$, and a decrease in $`N(0)`$. Calculating the first-order perturbative correction caused by the disorder shows that we must consider all these processes together. This is because the disorder-screened Coulomb interaction has a low-momentum singularity which leads to the separate effects being large; however, when they are added together this singularity is cancelled, and the actual effect is much smaller than might be naively expected. The final result has the form
$$\mathrm{ln}\left(\frac{T_c}{T_{c0}}\right)=\frac{1}{3}\frac{R_{\text{ }\text{ }\text{ }\text{ }\text{ }}}{R_0}\mathrm{ln}^3\left(\frac{1}{2\pi T_{c0}\tau }\right),$$
(2)
where $`R_0=2\pi h/e^2162k\mathrm{\Omega }`$ and $`\tau `$ is the elastic scattering time. We see that this curve is essentially “universal”, depending on only a single fitting parameter, $`\beta =\mathrm{ln}(1/2\pi T_{c0}\tau )`$. Experimentally it is found that $`T_c(R_{\text{ }\text{ }\text{ }\text{ }\text{ }})`$ curves from a wide variety of materials fit well to this equation, or extensions of it that allow for stronger disorder. The simplest extension simply consists of replacing $`T_{c0}`$ by $`T_c`$ on the right-hand side of Eq. (2), which leads to the cubic equation
$$x=\frac{t}{3}(\beta +x)^3,$$
(3)
where $`x=\mathrm{ln}(T_{c0}/T_c)`$ and $`t=R_{\text{ }\text{ }\text{ }\text{ }\text{ }}/R_0`$. This equation shows unphysical reentrance at strong disorder, an artefact which is removed by either a renormalization group treatment, or the use of a non-perturbative resummation technique to yield the formula
$$\mathrm{ln}\left(\frac{T_c}{T_{c0}}\right)=\frac{1}{\lambda }\frac{1}{2\sqrt{t}}\mathrm{ln}\left(\frac{1+\sqrt{t}/\lambda }{1\sqrt{t}/\lambda }\right).$$
(4)
The fact that most data can be fit to a single curve is pleasing in that it shows that the basic ingredients of our theory – disorder, BCS attraction and Coulomb repulsion – are correct. However it does not allow analysis of the sensitivity of the theory or experiment to changes in details of the system, such as the exact form of the phonon-mediated attraction. Moreover there are other theories which posit the importance of such details which give predictions that are equally in agreement with experiment. We see that the $`T_c(R_{\text{ }\text{ }\text{ }\text{ }\text{ }})`$ curves alone are not enough to allow consideration of the relative merits of different theories. What we would like to do is to add some additional parameter to the experimental system to give a whole new set of data – for example a family of $`T_c(R_{\text{ }\text{ }\text{ }\text{ }\text{ }})`$ curves for a single material as this new parameter is altered. Chervenak and Valles have recently performed an experiment of this type in which magnetic $`Gd`$ impurities are added to thin films of $`Pb_{0.9}Bi_{0.1}`$. This introduces the new feature of spin-flip scattering to the system, which is measured by the spin-flip scattering rate, $`1/\tau _s`$. The task of the theorist is to now make predictions for $`T_c`$ as a function of both $`R_{\text{ }\text{ }\text{ }\text{ }\text{ }}`$ (the measure of non-magnetic disorder), and $`1/\tau _s`$, (the measure of magnetic disorder), and to compare these to experiment.
In this paper we calculate the first-order perturbative correction to the transition temperature, $`T_c`$, of a superconductor with both non-magnetic and magnetic impurities. The model used consists of a featureless BCS attraction, $`\lambda `$, and a Coulomb repulsion, $`V_C(q)`$, between electrons which scatter off non-magnetic and magnetic impurities. The model is the simplest one that contains the essential physics, and its shortcoming of not considering the details of the attractive interaction is offset by the fact that we can consider all processes to a given order of perturbation theory. This is an important consideration in view of the cancellation of low-momentum singularities in the screened Coulomb potential discussed in the opening paragraph. In fact, an obvious question is whether this cancellation persists in the presence of magnetic impurities. We find that this is indeed the case, and so the details of the screened Coulomb interaction are removed, leading to a “universal” form for $`T_c(R_{\text{ }\text{ }\text{ }\text{ }\text{ }},1/\tau _s)`$.
The main result of the paper is that the pair-breaking rate per magnetic impurity, $`\alpha ^{}(R_{\text{ }\text{ }\text{ }\text{ }\text{ }})`$, defined by
$$\alpha ^{}(R_{\text{ }\text{ }\text{ }\text{ }\text{ }})=\frac{T_c(R_{\text{ }\text{ }\text{ }\text{ }\text{ }},0)T_c(R_{\text{ }\text{ }\text{ }\text{ }\text{ }},1/\tau _s)}{1/\tau _s}$$
(5)
is roughly independent of $`R_{\text{ }\text{ }\text{ }\text{ }\text{ }}`$ except near the superconductor-insulator transition, in agreement with experiment. This is confirmed both by first-order perturbation theory, and also by a non-perturbative resummation technique which we introduce to remove concerns about reentrance problems at stronger disorder. This agreement of the two theoretical approaches with each other and the experimental data gives us confidence in our results.
To calculate the correction to $`T_c`$ we use a collective mode formalism derived in a previous paper (which we refer to as I from now on) on the suppression of $`\mathrm{\Delta }`$ by non-magnetic disorder. The introduction of magnetic impurities means that we have to modify the formalism somewhat, and so we include most of the derivation in this paper. The method used in this paper to evaluate the correction to $`T_c`$ proceeds in three stages. First we find the first-order correction to the grand canonical potential, $`\mathrm{\Omega }_1(\mathrm{\Delta })`$, of the superconductor due to fluctuations of its collective modes. Then by minimizing the total grand canonical potential, $`\mathrm{\Omega }_0(\mathrm{\Delta })+\mathrm{\Omega }_1(\mathrm{\Delta })`$ with respect to the order parameter, $`\mathrm{\Delta }`$, we obtain the first-order correction to the order parameter self-consistency equation. Finally by setting $`\mathrm{\Delta }=0`$ we obtain the first-order correction to transition temperature $`T_c`$. The method we use has the advantage that it is impossible to “miss diagrams” since there is only one diagram in the $`\mathrm{\Omega }_1(\mathrm{\Delta })`$ calculation, and we also obtain the equation for $`\mathrm{\Delta }`$ at no extra cost. The equations for $`T_c`$ and $`\mathrm{\Delta }`$ must reduce to those of I when we set spin-flip scattering to zero, providing a useful consistency check. A key result of the calculation in I was that the singularity in the screened Coulomb potential persists below $`T_c`$, and that this singularity is cancelled in the formula for the suppression of $`\mathrm{\Delta }`$ in a similar manner to the cancellation in the formula for $`T_c`$. Therefore an important question is whether this singularity and cancellation remains when magnetic impurities are added. We show that this is indeed the case, and moreover that the cancellation is due to gauge invariance, and will occur in first-order perturbation theory in the presence of any kind of impurity scattering. In other words, it is not possible to obtain stronger suppression of the transition temperature by introducing some exotic scattering mechanism. It is also reassuring to know that an otherwise mysterious cancellation between diagrams has its physical origin in gauge invariance, and we hope that similar arguments may be applied to show that the result holds to all orders in perturbation theory.
The outline of the rest of the paper is as follows. In section II we derive the matrix formalism for superconductors with magnetic impurities, and the collective mode approach we will use. We derive the RPA screened bosonic propagators, and show that low-momentum singularities persist in the screened Coulomb propagator below $`T_c`$. In section III we derive the first order perturbative correction to the grand canonical potential $`\mathrm{\Omega }_1(\mathrm{\Delta })`$, and from this the correction to the order parameter $`\mathrm{\Delta }`$. In section IV we set $`\mathrm{\Delta }=0`$ to obtain the correction to transition temperature $`T_c`$. In section V we calculate $`T_c`$ numerically using both the perturbative results of section IV and a recently developed non-pertubative technique, and compare to experiment.
## II Derivation of the $`4\times 4`$ Matrix Formalism
Superconductivity with Magnetic Impurities:
We consider a system of electrons that scatter off static non-magnetic and magnetic impurities, and interact with each other via the long-range Coulomb interaction and the BCS attraction. The scattering from static impurities is described by Hamiltonian
$$H_{ei}=\underset{\alpha \beta }{}𝑑𝐱\psi _\alpha ^+(𝐱)\left\{\left[\frac{_x^2}{2m}+\underset{i}{}u_0(𝐱𝐱_𝐢)\right]\delta _{\alpha \beta }+\underset{j}{}J(𝐱𝐱_𝐣)𝐒_𝐣\sigma _{\alpha \beta }\right\}\psi _\beta (𝐱),$$
(6)
where $`\psi _\alpha ^+`$, $`\psi _\alpha `$ are the electron creation and annhilation operators, $`u_0(𝐱𝐱_𝐢)`$ is the impurity potential at $`𝐱`$ due to a non-magnetic impurity at $`𝐱_𝐢`$, $`𝐒_𝐣`$ is a magnetic impurity spin moment at $`𝐱_𝐣`$, and $`J(𝐱)`$ is the electron-impurity exchange coupling. The potential and spin-flip scattering rates are then given by
$`{\displaystyle \frac{1}{\tau _0}}`$ $`=`$ $`2\pi N(0)n_i|u_0|^2`$ (7)
$`{\displaystyle \frac{1}{\tau _s}}`$ $`=`$ $`2\pi N(0)n_jJ^2S(S+1),`$ (8)
where $`u_0`$ is the Fourier transform of $`u_0(𝐱)`$, and $`J`$ is the Fourier transform of $`J(𝐱)`$, both assumed independent of momentum, $`n_i`$ is the non-magnetic impurity density, and $`n_j`$ the magnetic impurity density.
The Coulomb repulsion between electrons is described by Hamiltonian
$$H_C=\underset{\alpha \beta }{}𝑑𝐱𝑑𝐱^{}\psi _\alpha ^+(𝐱)\psi _\alpha (𝐱)\frac{e^2}{|𝐱𝐱^{}|}\psi _\beta ^+(𝐱^{})\psi _\beta (𝐱^{}),$$
(9)
leading to a bare Coulomb propagator that is just the Fourier transform of the above potential.
The BCS attraction is described by the Hamiltonian
$$H_{BCS}=\lambda \underset{\alpha \beta }{}𝑑𝐱\psi _\alpha ^+(𝐱)\psi _\alpha (𝐱)\psi _\beta ^+(𝐱)\psi _\beta (𝐱),$$
(10)
corresponding to an instantaneous contact interaction $`\lambda \delta (𝐱𝐱^{})`$.
Having introduced the model Hamiltonian we need to describe the system, we discuss the standard four-dimensional matrix representation needed to describe a superconductor with magnetic impurities. We need four components to describe the two spin degrees of freedom, and the two types of correlation – the usual particle-hole correlation $`<\psi _{}\psi _{}^+>`$, amd the anomalous pairing correlation $`<\psi _{}\psi _{}>`$. We introduce the four-dimensional vector operator
$$\mathrm{\Psi }=\left(\begin{array}{c}\psi _{}\\ \psi _{}\\ \psi _{}^+\\ \psi _{}^+\end{array}\right);\mathrm{\Psi }^+=\left(\begin{array}{cccc}\psi _{}^+& \psi _{}^+& \psi _{}& \psi _{}\end{array}\right),$$
(11)
with matrix propagator
$$<\mathrm{\Psi }\mathrm{\Psi }^+>=\left[\begin{array}{cccc}<\psi _{}\psi _{}^+>& <\psi _{}\psi _{}^+>& <\psi _{}\psi _{}>& <\psi _{}\psi _{}>\\ <\psi _{}\psi _{}^+>& <\psi _{}\psi _{}^+>& <\psi _{}\psi _{}>& <\psi _{}\psi _{}>\\ <\psi _{}^+\psi _{}^+>& <\psi _{}^+\psi _{}^+>& <\psi _{}^+\psi _{}>& <\psi _{}^+\psi _{}>\\ <\psi _{}^+\psi _{}^+>& <\psi _{}^+\psi _{}^+>& <\psi _{}^+\psi _{}>& <\psi _{}^+\psi _{}>\end{array}\right].$$
(12)
In the normal state the temperature Green function is
$$G(k,i\omega )=\frac{1}{z\epsilon _k\tau _3\sigma _0},$$
(13)
where $`z=i\omega `$, $`\omega =(2n+1)\pi T`$ is a Fermi Matsubara frequency, and the $`\tau _i`$ and $`\sigma _i`$ are Pauli matrices operating on different spaces. The $`\sigma _i`$ operate in the usual spin space, whilst the $`\tau _i`$ operate in the Nambu (electron-hole) space. The diagrammatic rules are then the same as in the normal state except for the matrix structure of the Green function, and the presence of matrix $`\tau _3\sigma _0`$ at each interaction or impurity vertex due to the electron density operator being written in the form
$$\rho =\psi _{}^+\psi _{}+\psi _{}^+\psi _{}=\frac{1}{2}\mathrm{\Psi }^+\tau _3\sigma _0\mathrm{\Psi }.$$
(14)
The pairing correlations in the clean superconductor can then be taken into account self-consistently as shown in Fig. (1a). Making the ansatz $`\mathrm{\Sigma }=\mathrm{\Delta }\tau _1\sigma _3`$ for the self-energy, the Green function for the pure superconductor becomes
$$G_0(k,z)=\frac{1}{z\epsilon _k\tau _3\sigma _0\mathrm{\Delta }\tau _1\sigma _3}=\frac{z+\epsilon _k\tau _3\sigma _0+\mathrm{\Delta }\tau _1\sigma _3}{z^2\epsilon _k^2\mathrm{\Delta }^2},$$
(15)
and the diagram of Fig. (1a) gives self-energy
$`\mathrm{\Sigma }`$ $`=`$ $`\lambda T{\displaystyle \underset{\omega }{}}N(0){\displaystyle 𝑑\epsilon _k\frac{\tau _3\sigma _0(z+\epsilon _k\tau _3\sigma _0+\mathrm{\Delta }\tau _1\sigma _3)\tau _3\sigma _0}{\epsilon _k^2z^2+\mathrm{\Delta }^2}}`$ (16)
$`=`$ $`N(0)\lambda \mathrm{\Delta }\tau _1\sigma _3T{\displaystyle \underset{\omega }{}}{\displaystyle \frac{1}{\sqrt{\omega ^2+\mathrm{\Delta }^2}}},`$ (17)
which gives us the usual BCS self-consistency equation
$$1=N(0)\lambda T\underset{\omega }{}\frac{1}{\sqrt{\omega ^2+\mathrm{\Delta }^2}}.$$
(18)
We can treat the presence of non-magnetic and magnetic impurities by including an extra self-energy diagram to describe the dressing of the electron line by impurities as shown in Fig. (1b). We then make the ansatz that the pairing energy has the form $`\mathrm{\Sigma }_p=\mathrm{\Delta }\tau _1\sigma _3`$, and the impurity self-energy has the form $`\mathrm{\Sigma }_{imp}=(\overline{z}z)+(\overline{\mathrm{\Delta }}\mathrm{\Delta })\tau _1\sigma _3`$, so that the Green function for the dirty superconductor is
$$G_0(k,z)=\frac{\overline{z}+\epsilon _k\tau _3\sigma _0+\overline{\mathrm{\Delta }}\tau _1\sigma _3}{\overline{z}^2\epsilon _k^2\overline{\mathrm{\Delta }}^2},$$
(19)
which is just the Green function for the clean superconductor with $`z`$, $`\mathrm{\Delta }`$, replaced by $`\overline{z}`$, $`\overline{\mathrm{\Delta }}`$, respectively. Since the impurity line has the form
$$\mathrm{\Gamma }_0=\frac{1}{2\pi N(0)\tau _0}\tau _3\sigma _0\tau _3\sigma _0+\frac{1}{6\pi N(0)\tau _s}\left[\tau _0\sigma _1\tau _0\sigma _1+\tau _0\sigma _2\tau _0\sigma _2+\tau _3\sigma _3\tau _3\sigma _3\right],$$
(20)
we obtain the self-consistency equation for $`\overline{z}=i\overline{\omega }`$ and $`\overline{\mathrm{\Delta }}`$,
$$\overline{\omega }\omega =(\frac{1}{2\tau _0}+\frac{1}{2\tau _s})\frac{\overline{\omega }}{\sqrt{\overline{\omega }^2+\overline{\mathrm{\Delta }}^2}};\overline{\mathrm{\Delta }}\mathrm{\Delta }=(\frac{1}{2\tau _0}\frac{1}{2\tau _s})\frac{\overline{\mathrm{\Delta }}}{\sqrt{\overline{\omega }^2+\overline{\mathrm{\Delta }}^2}}.$$
(21)
The diagrammatic definition of the pairing energy, $`\mathrm{\Sigma }_p`$, leads to the same self-consistency equation for $`\mathrm{\Delta }`$ as in the pure case except that $`\omega `$, $`\mathrm{\Delta }`$, are replaced by $`\overline{\omega }`$, $`\overline{\mathrm{\Delta }}`$. In the absence of magnetic impurities – i.e. $`1/\tau _s=0`$ – we see that $`\overline{\omega }/\overline{\mathrm{\Delta }}=\omega /\mathrm{\Delta }`$, and the equation for $`\mathrm{\Delta }`$ is unchanged. This is Anderson’s theorem that superconductivity is unaffected by non-magnetic impurities at mean-field level. In the presence of magnetic impurities we see that $`\overline{\omega }/\overline{\mathrm{\Delta }}\omega /\mathrm{\Delta }`$, and if we define $`u=\overline{\omega }/\overline{\mathrm{\Delta }}`$, $`\zeta =1/\tau _s\mathrm{\Delta }`$, the problem reduces to solving the equation
$$\frac{\omega }{\mathrm{\Delta }}=u\left(1\zeta \frac{1}{\sqrt{u^2+1}}\right).$$
(22)
The self-consistency equation for $`\mathrm{\Delta }`$ then takes the form
$$1=N(0)\lambda T\underset{\omega }{}\frac{1}{\sqrt{u^2+1}},$$
(23)
and in particular if we set $`\mathrm{\Delta }=0`$ we get for $`T_c`$,
$$1=N(0)\lambda T_c\underset{\omega }{}\frac{1}{|\omega |+1/\tau _s}.$$
(24)
Subtracting off the equation for $`T_{c0}`$, the transition temperature in the absence of magnetic impurities leads to the famous result
$$\mathrm{log}\left(\frac{T_c}{T_{c0}}\right)=\psi \left(\frac{1}{2}\right)\psi \left(\frac{1}{2}+\frac{1}{2\pi T_c\tau _s}\right).$$
(25)
Collective Mode Formalism and RPA:
The idea of the collective mode formalism is to treat the screened interactions in the system as bosonic collective modes. The relevant bosonic operators are order parameter amplitude and phase, and electron density. These are the only modes that are coupled by a bare interaction – the BCS interaction for order parameter amplitude and phase, the Coulomb interaction for electron density. The main advantage of this approach is that we are able to treat order parameter fluctuations and the Coulomb interaction on an identical footing. This procedure may be formally carried out within the path integral theory of superconductivity by decoupling the two four-fermion interaction terms with the introduction of appropriate collective variables. This is discussed in detail in the paper of Eckern and Pelzer. The end result is that there are three effective bosonic modes, order parameter amplitude and phase and electronic density, and each can be written in the form
$$\widehat{O}_i=\frac{1}{2}\mathrm{\Psi }^+M_i\mathrm{\Psi },$$
(26)
with matrices $`M_\mathrm{\Delta }=\tau _1\sigma _3`$ for order parameter amplitude, $`M_\varphi =\tau _2\sigma _3`$ for order parameter phase, and $`M_\rho =\tau _3\sigma _0`$ for electronic density. Interactions occur by exchange of these collective modes, and so the effective interaction potential is now a $`3\times 3`$ matrix. The screened interaction is found from the equation
$$V_{ij}=V_{ij}^0+\underset{kl}{}V_{ik}^0\mathrm{\Pi }_{kl}V_{lj},$$
(27)
as shown in Fig. (2). Here $`\mathrm{\Pi }_{ij}`$ is the polarization operator and $`V_{ij}^0`$ is the bare interaction matrix which is given by
$$V^0=\left(\begin{array}{ccc}\lambda /2& 0& 0\\ 0& \lambda /2& 0\\ 0& 0& 4\pi e^2/q^2\end{array}\right),$$
(28)
the BCS attraction being split equally between the two order parameter modes. The only new diagrammatic feature is that any of the matrices $`M_i`$ can now appear at an interaction vertex, corresponding to interaction with the collective variable described by that matrix.
In order to carry out the calculations in section III, we need the impurity dressed RPA polarization bubbles, $`\mathrm{\Pi }_{ij}`$, shown in Fig. (3). To evaluate the polarization bubble $`\mathrm{\Pi }_{ij}`$ we must first evaluate the geometric series
$$\mathrm{\Pi }=S+S\mathrm{\Gamma }_0S+S\mathrm{\Gamma }_0S\mathrm{\Gamma }_0S+\mathrm{},$$
(29)
where $`\mathrm{\Gamma }_0`$ is the impurity line defined in Eqn. (20), and $`S`$ is the momentum sum of a direct product of Green functions
$$S=\underset{k}{}G(k,i\omega )G(k+q,i\omega +i\mathrm{\Omega }).$$
(30)
Since we do not need the complete matrix structure of $`\mathrm{\Pi }`$, but just its traces with two matrices from the set $`\tau _1\sigma _3`$, $`\tau _2\sigma _3`$, $`\tau _3\sigma _0`$, we actually evaluate the impurity dressed vertices $`\mathrm{\Pi }_j`$ which have one matrix from the above set inserted between the two terms of the direct product in $`\mathrm{\Pi }`$. These satisfy the equations
$$\mathrm{\Pi }_j=S_j+S\mathrm{\Gamma }_0\mathrm{\Pi }_j,$$
(31)
where $`S_j`$ is obtained by inserting the matrix $`M_j`$ between the two terms of the direct product in $`S`$. These equations are solved in Appendix A to obtain the results
$`\mathrm{\Pi }_\mathrm{\Delta }`$ $`=`$ $`{\displaystyle \frac{\pi N(0)}{D_+}}\left[{\displaystyle \frac{UU^{}+uu^{}1}{UU^{}}}{\displaystyle \frac{i(u^{}+u)}{UU^{}}}\tau _1\sigma _3\right]\tau _1\sigma _3`$ (32)
$`\mathrm{\Pi }_\varphi `$ $`=`$ $`{\displaystyle \frac{\pi N(0)}{D_{}}}\left[{\displaystyle \frac{UU^{}+uu^{}+1}{UU^{}}}{\displaystyle \frac{i(u^{}u)}{UU^{}}}\tau _1\sigma _3\right]\tau _2\sigma _3`$ (33)
$`\mathrm{\Pi }_\varphi `$ $`=`$ $`{\displaystyle \frac{\pi N(0)}{D_{}}}\left[{\displaystyle \frac{UU^{}uu^{}1}{UU^{}}}+{\displaystyle \frac{i(u^{}u)}{UU^{}}}\tau _1\sigma _3\right]\tau _3\sigma _0,`$ (34)
where $`U=\sqrt{u^2+1}`$, $`u^{}=u(\omega +\mathrm{\Omega })=u(\omega ^{})`$, $`U^{}=\sqrt{u^2+1}`$ and
$$D_\pm =\left[Dq^2+\mathrm{\Delta }U+\mathrm{\Delta }U^{}+\frac{1}{\tau _s}\left(\frac{uu^{}1}{UU^{}}1\right)\right].$$
(35)
We can finally obtain the non-zero polarization bubbles $`\mathrm{\Pi }_{ij}`$ by inserting the second matrix from the set $`\tau _1\sigma _3`$, $`\tau _2\sigma _3`$, $`\tau _3\sigma _0`$ into $`\mathrm{\Pi }_j`$, taking the trace, and recalling the factor $`1`$ for a fermion loop. This yields
$`\mathrm{\Pi }_{\mathrm{\Delta }\mathrm{\Delta }}(q,\mathrm{\Omega })`$ $`=`$ $`\pi N(0)T{\displaystyle \underset{\omega }{}}\left[{\displaystyle \frac{UU^{}+uu^{}1}{UU^{}}}\right]{\displaystyle \frac{1}{\left(Dq^2+\mathrm{\Delta }U+\mathrm{\Delta }U^{}{\displaystyle \frac{1}{\tau _s}}\left[{\displaystyle \frac{UU^{}uu^{}+1}{UU^{}}}\right]\right)}}`$ (36)
$`\mathrm{\Pi }_{\varphi \varphi }(q,\mathrm{\Omega })`$ $`=`$ $`\pi N(0)T{\displaystyle \underset{\omega }{}}\left[{\displaystyle \frac{UU^{}+uu^{}+1}{UU^{}}}\right]{\displaystyle \frac{1}{\left(Dq^2+\mathrm{\Delta }U+\mathrm{\Delta }U^{}{\displaystyle \frac{1}{\tau _s}}\left[{\displaystyle \frac{UU^{}uu^{}1}{UU^{}}}\right]\right)}}`$ (37)
$`\mathrm{\Pi }_{\rho \rho }(q,\mathrm{\Omega })`$ $`=`$ $`\pi N(0)T{\displaystyle \underset{\omega }{}}\left[{\displaystyle \frac{UU^{}uu^{}1}{UU^{}}}\right]{\displaystyle \frac{1}{\left(Dq^2+\mathrm{\Delta }U+\mathrm{\Delta }U^{}{\displaystyle \frac{1}{\tau _s}}\left[{\displaystyle \frac{UU^{}uu^{}1}{UU^{}}}\right]\right)}}+N(0)`$ (38)
$`\mathrm{\Pi }_{\varphi \rho }(q,\mathrm{\Omega })`$ $`=`$ $`\pi N(0)T{\displaystyle \underset{\omega }{}}\left[{\displaystyle \frac{u^{}u}{UU^{}}}\right]{\displaystyle \frac{1}{\left(Dq^2+\mathrm{\Delta }U+\mathrm{\Delta }U^{}{\displaystyle \frac{1}{\tau _s}}\left[{\displaystyle \frac{UU^{}uu^{}1}{UU^{}}}\right]\right)}}=\mathrm{\Pi }_{\rho \varphi }(q,\mathrm{\Omega }).`$ (39)
We note that if we set $`1/\tau _s=0`$ in Eqn. (36) we will obtain exactly the results found in I, as of course we must. The screened potentials $`V_{ij}`$ are then given by
$$V=\left[\begin{array}{ccc}(\lambda ^1+\mathrm{\Pi }_{\mathrm{\Delta }\mathrm{\Delta }})^1& 0& 0\\ 0& [(2V_C(q))^1+\mathrm{\Pi }_{\rho \rho }]/𝒟& \mathrm{\Pi }_{\varphi \rho }/𝒟\\ 0& \mathrm{\Pi }_{\varphi \rho }/𝒟& (\lambda ^1+\mathrm{\Pi }_{\varphi \varphi })/𝒟\end{array}\right],$$
(40)
where
$$𝒟(\lambda ^1+\mathrm{\Pi }_{\varphi \varphi })[(2V_C(q))^1+\mathrm{\Pi }_{\rho \rho }]+\mathrm{\Pi }_{\varphi \rho }^2.$$
(41)
The coupling between phase and density fluctuations caused by the non-zero value of $`\mathrm{\Pi }_{\varphi \rho }=\mathrm{\Pi }_{\rho \varphi }`$ is a manifestation of gauge invariance.
We can next show that the propagators $`V_{\varphi \varphi }`$, $`V_{\varphi \rho }`$ and $`V_{\rho \rho }`$ all have a low-momentum singularity of the form $`1/q^{d1}`$ for all non-zero frequencies and all temperatures. In other words, these propagators have the same long-range behavior as the unscreened Coulomb potential. An analogous situation occurs for the disorder screened potential in the normal metal, where the singularity is known to strongly affect the properties of the system. To show the existence of the singularity we simply need to show that the denominator $`𝒟`$ vanishes at $`q=0`$ for all $`\mathrm{\Omega }0`$ and all $`T`$. Since $`V_C(q)1/q^{d1}`$, we need only prove that
$$[\lambda ^1\mathrm{\Pi }_{\varphi \varphi }(0,\mathrm{\Omega })]\mathrm{\Pi }_{\rho \rho }(0,\mathrm{\Omega })+\mathrm{\Pi }_{\varphi \rho }(0,\mathrm{\Omega })^2=0.$$
(42)
This is proved in Appendix B where we show that
$$[\lambda ^1+\mathrm{\Pi }_{\varphi \varphi }(0,\mathrm{\Omega })]=\frac{\mathrm{\Omega }}{2\mathrm{\Delta }}\mathrm{\Pi }_{\varphi \rho }(0,\mathrm{\Omega }),\mathrm{\Pi }_{\varphi \rho }(0,\mathrm{\Omega })=\frac{\mathrm{\Omega }}{\mathrm{\Delta }}\mathrm{\Pi }_{\rho \rho }(0,\mathrm{\Omega }),$$
(43)
by direct calculation. We also show that this result is guaranteed by gauge invariance and is therefore true no matter which scattering mechanisms we include.
## III First Order Correction to Grand Potential and Order-Parameter Self-Consistency Equation
In this section we evaluate the first-order perturbation correction to the grand potential, $`\mathrm{\Omega }_1(\mathrm{\Delta })`$. By minimising the sum $`\mathrm{\Omega }_0(\mathrm{\Delta })+\mathrm{\Omega }_1(\mathrm{\Delta })`$ with respect to $`\mathrm{\Delta }`$ we obtain the corresponding correction to the order parameter self-consistency equation. This method was first used for the system with only non-magnetic impurities by Eckern and Pelzer, and we choose to use it as it involves the smallest number of diagrams. The same result for $`T_c`$ can also be obtained using the Eliashberg diagrams for $`\mathrm{\Delta }`$ shown in Fig. (5), or the pair propagator diagrams shown in Fig. (6).
The diagram for the first order correction to the grand potential simply consists of the “string of bubbles” diagram shown in Fig. (4) Since the polarization bubbles in this diagram are just those evaluated in the previous section, we have all the information we need to derive $`\mathrm{\Omega }_1(\mathrm{\Delta })`$. The only thing we need to remember is the extra symmetry factor of $`1/n`$ required for the diagram with $`n`$ bubbles. So whereas previously the RPA equation involved the series
$$V=V_0+V_0\mathrm{\Pi }V_0+V_0\mathrm{\Pi }V_0\mathrm{\Pi }V_0+\mathrm{}=(V_0^1\mathrm{\Pi })^1,$$
(44)
it now becomes
$$\mathrm{\Omega }_1=V_0\mathrm{\Pi }+\frac{1}{2}V_0\mathrm{\Pi }V_0\mathrm{\Pi }+\frac{1}{3}V_0\mathrm{\Pi }V_0\mathrm{\Pi }V_0\mathrm{\Pi }+\mathrm{}=\mathrm{log}[V_0^1\mathrm{\Pi }]+\mathrm{log}[V_0^1].$$
(45)
After summing over all the internal variables – momentum $`q`$, Bose Matsubara frequency $`\mathrm{\Omega }`$ and the three bosonic modes – we end up with the final expression for $`\mathrm{\Omega }_1`$,
$$\mathrm{\Omega }_1=T\underset{\mathrm{\Omega }}{}\underset{q}{}\left\{\mathrm{log}(\lambda ^1+\mathrm{\Pi }_{\mathrm{\Delta }\mathrm{\Delta }}(q,\mathrm{\Omega }))+\mathrm{log}\left[(\lambda ^1+\mathrm{\Pi }_{\varphi \varphi }(q,\mathrm{\Omega }))[(2V_C(q))^1+\mathrm{\Pi }_{\rho \rho }(q,\mathrm{\Omega })]+\mathrm{\Pi }_{\varphi \rho }(q,\mathrm{\Omega })^2\right]\right\}.$$
(46)
To proceed we need to minimise the total grand potential,
$$\frac{}{\mathrm{\Delta }}\left\{\mathrm{\Omega }_0(\mathrm{\Delta })+\mathrm{\Omega }_1(\mathrm{\Delta })\right\}=0.$$
(47)
The mean-field grand potential, $`\mathrm{\Omega }_0(\mathrm{\Delta })`$, is given by
$$\mathrm{\Omega }_0(\mathrm{\Delta })=N(0)\frac{\mathrm{\Delta }^2}{\lambda }+T\underset{\omega }{}\underset{k}{}\text{Tr}\left[\mathrm{log}(\overline{z}\xi _k\tau _3\sigma _0\overline{\mathrm{\Delta }}\tau _1\sigma _3)\right],$$
(48)
and after taking the derivative of $`\mathrm{\Omega }_0(\mathrm{\Delta })`$ with respect to $`\mathrm{\Delta }`$, we see that Eqn. (47) takes the form
$$\frac{1}{\lambda }\pi N(0)T\underset{\omega }{}\frac{1}{U}=\frac{1}{\mathrm{\Delta }}\frac{\mathrm{\Omega }_1}{\mathrm{\Delta }}.$$
(49)
The next step in the procedure is to evaluate $`\mathrm{\Omega }_1(\mathrm{\Delta })/\mathrm{\Delta }`$. From Eqn. (46) we see that $`/\mathrm{\Delta }`$ will be acting upon the $`\mathrm{\Pi }_{ij}`$ to give
$$\frac{\mathrm{\Omega }_1}{\mathrm{\Delta }}=T\underset{\mathrm{\Omega }}{}\underset{q}{}\left\{\frac{\mathrm{\Pi }_{\mathrm{\Delta }\mathrm{\Delta }}}{\mathrm{\Delta }}V_{\mathrm{\Delta }\mathrm{\Delta }}+\frac{\mathrm{\Pi }_{\varphi \varphi }}{\mathrm{\Delta }}V_{\varphi \varphi }2\frac{\mathrm{\Pi }_{\varphi \rho }}{\mathrm{\Delta }}V_{\varphi \rho }+\frac{\mathrm{\Pi }_{\rho \rho }}{\mathrm{\Delta }}V_{\rho \rho }\right\}.$$
(50)
From Eqn. (36) we see that the $`/\mathrm{\Delta }`$ can act either on the coherence factor or on the denominator in the expression for $`\mathrm{\Pi }_{ij}`$. Acting on the denominator gives a result proportional to the denominator squared, corresponding to the two-ladder diagrams of Fig. (5a) in the Eliashberg approach. Similarly acting on the coherence factor leads to the one-ladder diagrams of Fig. (5b). We note that the explicit evaluation of the three-ladder diagrams of Fig. (5c) will give a zero result.
The only difficulty in taking the derivatives of the polarization bubbles, $`\mathrm{\Pi }_{ij}`$, with respect to $`\mathrm{\Delta }`$ is that the quantity $`u(\omega )`$ present in all these equations satisfies the transcendental Eqn. (22). In Appendix C we evaluate the derivatives of the $`\mathrm{\Pi }_{ij}`$ to obtain the results
$`{\displaystyle \frac{\mathrm{\Pi }_{\mathrm{\Delta }\mathrm{\Delta }}}{\mathrm{\Delta }}}`$ $`=`$ $`2\pi N(0)T{\displaystyle \underset{\omega }{}}\left\{\left[1{\displaystyle \frac{\zeta }{U^3}}\right]^1\times {\displaystyle \frac{1}{\mathrm{\Delta }}}{\displaystyle \frac{u(u^{}+u)}{U^3U^{}}}{\displaystyle \frac{1}{D_+}}+\left(1+{\displaystyle \frac{uu^{}1}{UU^{}}}\right)\left\{{\displaystyle \frac{1}{U}}{\displaystyle \frac{\zeta }{U^2}}\left(1+{\displaystyle \frac{u(u^{}u)}{UU^{}}}\right)\right\}{\displaystyle \frac{1}{D_+^2}}\right\}`$ (51)
$`{\displaystyle \frac{\mathrm{\Pi }_{\varphi \varphi }}{\mathrm{\Delta }}}`$ $`=`$ $`2\pi N(0)T{\displaystyle \underset{\omega }{}}\left\{\left[1{\displaystyle \frac{\zeta }{U^3}}\right]^1\times {\displaystyle \frac{1}{\mathrm{\Delta }}}{\displaystyle \frac{u(u^{}u)}{U^3U^{}}}{\displaystyle \frac{1}{D_{}}}+\left(1+{\displaystyle \frac{uu^{}+1}{UU^{}}}\right)\left\{{\displaystyle \frac{1}{U}}{\displaystyle \frac{\zeta }{U^2}}\left(1+{\displaystyle \frac{u(u^{}+u)}{UU^{}}}\right)\right\}{\displaystyle \frac{1}{D_{}^2}}\right\}`$ (52)
$`{\displaystyle \frac{\mathrm{\Pi }_{\rho \rho }}{\mathrm{\Delta }}}`$ $`=`$ $`2\pi N(0)T{\displaystyle \underset{\omega }{}}\left\{\left[1{\displaystyle \frac{\zeta }{U^3}}\right]^1\times {\displaystyle \frac{1}{\mathrm{\Delta }}}{\displaystyle \frac{u(u^{}u)}{U^3U^{}}}{\displaystyle \frac{1}{D_{}}}\left(1{\displaystyle \frac{uu^{}+1}{UU^{}}}\right)\left\{{\displaystyle \frac{1}{U}}{\displaystyle \frac{\zeta }{U^2}}\left(1+{\displaystyle \frac{u(u^{}+u)}{UU^{}}}\right)\right\}{\displaystyle \frac{1}{D_{}^2}}\right\}`$ (53)
$`{\displaystyle \frac{\mathrm{\Pi }_{\varphi \rho }}{\mathrm{\Delta }}}`$ $`=`$ $`2\pi N(0)T{\displaystyle \underset{\omega }{}}\left\{\left[1{\displaystyle \frac{\zeta }{U^3}}\right]^1\times {\displaystyle \frac{1}{\mathrm{\Delta }}}{\displaystyle \frac{u(uu^{}+1)}{U^3U^{}}}{\displaystyle \frac{1}{D_{}}}\left({\displaystyle \frac{u^{}u}{UU^{}}}\right)\left\{{\displaystyle \frac{1}{U}}{\displaystyle \frac{\zeta }{U^2}}\left(1+{\displaystyle \frac{u(u^{}+u)}{UU^{}}}\right)\right\}{\displaystyle \frac{1}{D_{}^2}}\right\}.`$ (54)
These formulas together with Eqn. (50) lead to our final result for the first order correction to the order parameter self-consistency equation:
$`{\displaystyle \frac{1}{N(0)\lambda }}\pi T{\displaystyle \underset{\omega }{}}{\displaystyle \frac{1}{U}}={\displaystyle \frac{\pi }{2}}T{\displaystyle \underset{\omega }{}}T{\displaystyle \underset{\mathrm{\Omega }}{}}{\displaystyle \underset{q}{}}[1{\displaystyle \frac{\zeta }{U^3}}]^1\times `$ (55)
$`\{[{\displaystyle \frac{1}{\mathrm{\Delta }^2}}{\displaystyle \frac{1}{D_+}}{\displaystyle \frac{u(u^{}+u)}{U^3U^{}}}+{\displaystyle \frac{1}{\mathrm{\Delta }}}{\displaystyle \frac{1}{D_+^2}}{\displaystyle \frac{1}{U}}\{1{\displaystyle \frac{\zeta }{U}}(1+{\displaystyle \frac{u(u^{}u)}{UU^{}}})\}(1+{\displaystyle \frac{uu^{}1}{UU^{}}})]V_{\mathrm{\Delta }\mathrm{\Delta }}(q,\mathrm{\Omega })`$ (56)
$`+{\displaystyle \frac{1}{\mathrm{\Delta }^2}}{\displaystyle \frac{1}{D_{}}}\left[{\displaystyle \frac{u(u^{}u)}{U^3U^{}}}V_{\varphi \varphi }(q,\mathrm{\Omega }){\displaystyle \frac{2u(uu^{}+1)}{U^3U^{}}}V_{\varphi \rho }(q,\mathrm{\Omega })+{\displaystyle \frac{u(u^{}u)}{U^3U^{}}}V_{\rho \rho }(q,\mathrm{\Omega })\right]`$ (57)
$`+{\displaystyle \frac{1}{\mathrm{\Delta }}}{\displaystyle \frac{1}{D_{}^2}}\{1{\displaystyle \frac{\zeta }{U}}(1+{\displaystyle \frac{u(u^{}u)}{UU^{}}})\}[(1+{\displaystyle \frac{uu^{}+1}{UU^{}}})V_{\varphi \varphi }(q,\mathrm{\Omega })+{\displaystyle \frac{2(u^{}u)}{UU^{}}}V_{\varphi \rho }(q,\mathrm{\Omega })(1{\displaystyle \frac{uu^{}+1}{UU^{}}})V_{\rho \rho }]\}.`$ (58)
The above formula is valid for all temperatures $`0TT_c`$, but we are usually interested in the special cases $`T=0`$ and $`\mathrm{\Delta }=0`$ (i.e. $`T=T_c`$). In these two cases the sum over $`\omega `$ on the LHS can be performed analytically to yield the two simple forms
$`\mathrm{log}\left({\displaystyle \frac{\mathrm{\Delta }(0)}{\mathrm{\Delta }_0(0)}}\right)`$ $`=`$ $`{\displaystyle \frac{\pi }{2}}T{\displaystyle \underset{\mathrm{\Omega }}{}}{\displaystyle \underset{q}{}}T{\displaystyle \underset{\omega }{}}\mathrm{}`$ (59)
$`\mathrm{log}\left({\displaystyle \frac{T_c}{T_{c0}}}\right)`$ $`=`$ $`{\displaystyle \frac{\pi }{2}}T{\displaystyle \underset{\mathrm{\Omega }}{}}{\displaystyle \underset{q}{}}T{\displaystyle \underset{\omega }{}}\mathrm{}`$ (60)
Having noted the presence of the $`1/q^{d1}`$ singularities in the potentials $`V_{\varphi \varphi }`$, $`V_{\varphi \rho }`$ and $`V_{\rho \rho }`$, we should now see whether the terms in Eqn. (55) containing these singularities cancel out. If we go back to Eqn. (46) for the correction to the grand potential, we see that the term $`𝒟`$ that goes as $`q^{d1}`$ is inside a logarithm. Since $`q^{d1}`$ occurs as a product, we can simply take off the term $`\mathrm{log}(q^{d1})`$, and upon differentiating with respect to $`\mathrm{\Delta }`$ should get zero. In other words we naively expect no singular term in Eqn. (55). However this is not quite correct since to prove that the denominator $`𝒟`$ vanished at $`q=0`$ we needed to replace $`\lambda ^1`$ using the mean-field self-consistency equation, Eqn. (23). We note that although the two sides of Eqn. (23) are numerically equal in the mean-field case, their dependences on $`\mathrm{\Delta }`$ differ – $`\lambda ^1`$ gives zero under $`/\mathrm{\Delta }`$, whilst $`1/U`$ does not. This discrepancy then leads to the only singular term in Eqn. (55), which may be written
$$\mathrm{ln}\left(\frac{\mathrm{\Delta }}{\mathrm{\Delta }_0}\right)_{mf}=\frac{1}{4}\pi T\underset{\omega }{}\frac{1}{\mathrm{\Delta }^2U^3}[1\frac{\zeta }{U^3}]^1T\underset{\mathrm{\Omega }}{}\underset{q}{}V_{\varphi \varphi }(q,\mathrm{\Omega }).$$
(61)
Since this term tends to half the pair propagator contribution to the suppression of $`T_c`$ when we let $`\mathrm{\Delta }0`$, we interpret it as the phase fluctuation contribution. It is singular because of the Mermin-Wagner-Hohenberg theorem which tells us that we cannot have broken symmetry states in 2D systems at finite temperature. In the following we will be mainly interested in the correction to $`T_c`$ due to Coulomb interaction and so will ignore this term.
## IV First Order Correction to the Transition Temperature
We can now evaluate the first order correction to the transition temperature by linearizing the order parameter self-consistency equation with respect to $`\mathrm{\Delta }`$. The former can also be obtained directly from the normal state by calculating the pair propagator $`L(q,\mathrm{\Omega })`$ to first order, and looking for the instability at $`q=\mathrm{\Omega }=0`$. $`L`$ is given by
$$L^1(q,\mathrm{\Omega })=\lambda ^1+P(q,\mathrm{\Omega }),$$
(62)
where $`P(q,\mathrm{\Omega })`$ is the pair polarization bubble. The zeroth order polarization bubble $`P_0(q,\mathrm{\Omega })`$ is shown in Fig. (6b) and leads to the mean-field result
$`L_0^1(q,\mathrm{\Omega })`$ $`=`$ $`N(0)\left[\mathrm{log}\left({\displaystyle \frac{T}{T_{c00}}}\right)+\psi \left({\displaystyle \frac{1}{2}}+{\displaystyle \frac{1}{2\pi T\tau _s}}+{\displaystyle \frac{Dq^2+|\mathrm{\Omega }|}{4\pi T}}\right)\psi \left({\displaystyle \frac{1}{2}}\right)\right]`$ (63)
$`=`$ $`N(0)\left[\mathrm{log}\left({\displaystyle \frac{T}{T_{c0}}}\right)+\psi \left({\displaystyle \frac{1}{2}}+{\displaystyle \frac{1}{2\pi T\tau _s}}+{\displaystyle \frac{Dq^2+|\mathrm{\Omega }|}{4\pi T}}\right)\psi \left({\displaystyle \frac{1}{2}}+{\displaystyle \frac{1}{2\pi T\tau _s}}\right)\right],`$ (64)
where $`T_{c00}`$ is the BCS transition temperature (the mean field value in the absence of magnetic impurities), and $`T_{c0}`$ is the mean field value for the system with magnetic impurities. A correction to the polarization operator $`\delta P(0,0)`$ will lead to a change in the transition temperature, which is defined as the temperature at which the denominator of $`L`$ becomes zero, given by
$$\mathrm{log}\left(\frac{T_c}{T_{c0}}\right)=\frac{\delta P(0,0)}{N(0)}.$$
(65)
If we look at Fig. (6) we see that there are 7 diagrams which contribute to the first order correction to $`T_c`$. We will set $`\mathrm{\Delta }=0`$ in the order parameter result of Eqn. (55) to get the transition temperature equation, and we will be able to identify the contribution that comes from each of the $`P_i`$ diagrams. When we set $`\mathrm{\Delta }0`$, then $`\mathrm{\Delta }u(|\omega |+1/\tau _s)\text{sgn}(\omega )`$; $`\mathrm{\Delta }U|\omega |+1/\tau _s`$; $`V_{\rho \rho }2V_C`$; $`V_{\mathrm{\Delta }\mathrm{\Delta }}`$ and $`V_{\varphi \varphi }L`$; $`V_{\varphi \rho }2\mathrm{\Pi }_{\varphi \rho }LV_C`$. The coherence factors then become Heaviside functions that set the relative signs of the frequencies
$`1{\displaystyle \frac{uu^{}\pm 1}{UU^{}}}`$ $``$ $`2\theta (\omega (\omega +\mathrm{\Omega }))`$ (66)
$`1+{\displaystyle \frac{uu^{}\pm 1}{UU^{}}}`$ $``$ $`2\theta (\omega (\omega +\mathrm{\Omega })).`$ (67)
The two denominators, $`D_\pm `$, both become $`Dq^2+|\mathrm{\Omega }|`$ for $`\omega `$, $`\omega +\mathrm{\Omega }`$ of opposite sign; $`Dq^2+|2\omega +\mathrm{\Omega }|+2/\tau _s`$ for $`\omega `$, $`\omega +\mathrm{\Omega }`$ of the same sign. Making all these substitutions leads to
$`P_1`$ $`=`$ $`\pi N(0)T{\displaystyle \underset{\omega }{}}T{\displaystyle \underset{\mathrm{\Omega }}{}}{\displaystyle \underset{q}{}}\left[{\displaystyle \frac{1}{(|\omega |+1/\tau _s)^2}}{\displaystyle \frac{1}{(Dq^2+|\mathrm{\Omega }|)}}+{\displaystyle \frac{2}{(|\omega |+1/\tau _s)}}{\displaystyle \frac{1}{(Dq^2+|\mathrm{\Omega }|)^2}}\right]V_C(q,\mathrm{\Omega })\theta (\omega (\omega +\mathrm{\Omega }))`$ (68)
$`P_2`$ $`=`$ $`\pi N(0)T{\displaystyle \underset{\omega }{}}T{\displaystyle \underset{\mathrm{\Omega }}{}}{\displaystyle \underset{q}{}}{\displaystyle \frac{1}{(|\omega |+1/\tau _s)^2}}{\displaystyle \frac{1}{Dq^2+|2\omega +\mathrm{\Omega }|+2/\tau _s}}V_C(q,\mathrm{\Omega })\theta (\omega (\omega +\mathrm{\Omega }))`$ (69)
$`P_3`$ $`=`$ $`\pi N(0)T{\displaystyle \underset{\omega }{}}T{\displaystyle \underset{\mathrm{\Omega }}{}}{\displaystyle \underset{q}{}}{\displaystyle \frac{1}{(|\omega |+1/\tau _s)(|\omega +\mathrm{\Omega }|+1/\tau _s)}}\left[{\displaystyle \frac{1}{(Dq^2+|\mathrm{\Omega }|)}}+{\displaystyle \frac{2/\tau _s}{(Dq^2+|\mathrm{\Omega }|)^2}}\right]V_C(q,\mathrm{\Omega })\theta (\omega (\omega +\mathrm{\Omega }))`$ (70)
$`P_4`$ $`=`$ $`\pi N(0)T{\displaystyle \underset{\omega }{}}T{\displaystyle \underset{\mathrm{\Omega }}{}}{\displaystyle \underset{q}{}}{\displaystyle \frac{1}{(|\omega |+1/\tau _s)(|\omega +\mathrm{\Omega }|+1/\tau _s)}}{\displaystyle \frac{1}{(Dq^2+|2\omega +\mathrm{\Omega }|+2/\tau _s)}}V_C(q,\mathrm{\Omega })\theta (\omega (\omega +\mathrm{\Omega }))`$ (71)
$`P_5`$ $`=`$ $`\pi N(0)T{\displaystyle \underset{\omega }{}}T{\displaystyle \underset{\mathrm{\Omega }}{}}{\displaystyle \underset{q}{}}{\displaystyle \frac{1}{(|\omega |+1/\tau _s)^2}}\left[{\displaystyle \frac{1}{(Dq^2+|2\omega +\mathrm{\Omega }|+2/\tau _s)}}+{\displaystyle \frac{|\omega |1/\tau _s}{(Dq^2+|2\omega +\mathrm{\Omega }|+2/\tau _s)^2}}\right]L(q,\mathrm{\Omega })\theta (\omega (\omega +\mathrm{\Omega }))`$ (72)
$`P_6`$ $`=`$ $`\pi N(0)T{\displaystyle \underset{\omega }{}}T{\displaystyle \underset{\mathrm{\Omega }}{}}{\displaystyle \underset{q}{}}{\displaystyle \frac{1}{(|\omega |+1/\tau _s)^2(Dq^2+|\mathrm{\Omega }|)}}L(q,\mathrm{\Omega })\theta (\omega (\omega +\mathrm{\Omega }))`$ (73)
$`P_7`$ $`=`$ $`4\pi ^2N(0)^2T{\displaystyle \underset{\mathrm{\Omega }}{}}{\displaystyle \underset{q}{}}\left[T{\displaystyle \underset{\omega }{}}{\displaystyle \frac{\text{sgn}(\omega +\mathrm{\Omega })}{(|\omega |+1/\tau _s)(Dq^2+|\omega |+|\omega +\mathrm{\Omega }|+2/\tau _s)\theta (\omega (\omega +\mathrm{\Omega })))}}\right]^2V_C(q,\mathrm{\Omega })L(q,\mathrm{\Omega }).`$ (74)
The assignment of terms to the polarization bubble diagram they would arise from if we had done the calculation by that method is unique, and can be summarised below:
* $`P_1`$ : term proportional to $`V_C`$, $`\theta (\omega (\omega +\mathrm{\Omega }))`$ with no $`|\omega +\mathrm{\Omega }|+1/\tau _s`$ denominator.
* $`P_2`$ : term proportional to $`V_C`$, $`\theta (\omega (\omega +\mathrm{\Omega }))`$ with no $`|\omega +\mathrm{\Omega }|+1/\tau _s`$ denominator.
* $`P_3`$ : term proportional to $`V_C`$, $`\theta (\omega (\omega +\mathrm{\Omega }))`$ with $`|\omega +\mathrm{\Omega }|+1/\tau _s`$ denominator.
* $`P_4`$ : term proportional to $`V_C`$, $`\theta (\omega (\omega +\mathrm{\Omega }))`$ with no $`|\omega +\mathrm{\Omega }|+1/\tau _s`$ denominator.
* $`P_5`$ : term proportional to $`L`$ and $`\theta (\omega (\omega +\mathrm{\Omega }))`$.
* $`P_6`$ : term proportional to $`L`$ and $`\theta (\omega (\omega +\mathrm{\Omega }))`$.
* $`P_7`$ : term proportional to $`LV_C`$.
We find that these reduce to the results of I when we set $`1/\tau _s0`$, providing a useful consistency check on the present calculation.
To evaluate the correction to $`T_c`$ we split it into two parts: the Coulomb part consisting of those parts that contain a Coulomb propagator, ($`P_1P_4`$, $`P_7`$), and consequently require special attention at $`q=0`$, and the fluctuation part consisting of those terms that contain only a fluctuation propagator, ($`P_5`$, $`P_6`$). Performing the $`\omega `$-sum first we get for the Coulomb part
$`\mathrm{log}\left({\displaystyle \frac{T_c}{T_{c0}}}\right)=T_c{\displaystyle \underset{\mathrm{\Omega }}{}}{\displaystyle \underset{q}{}}\{{\displaystyle \frac{1}{2\pi T_c}}{\displaystyle \frac{Dq^2}{\mathrm{\Omega }^2(Dq^2)^2}}\psi ^{}({\displaystyle \frac{1}{2}}+{\displaystyle \frac{1}{2\pi T_c\tau _s}}+{\displaystyle \frac{|\mathrm{\Omega }|}{2\pi T_c}})`$ (75)
$`+`$ $`\left[{\displaystyle \frac{2Dq^2[\mathrm{\Omega }^2+(Dq^2)^2]}{|\mathrm{\Omega }|[\mathrm{\Omega }^2(Dq^2)^2]^2}}{\displaystyle \frac{\left({\displaystyle \frac{2}{\tau _s}}\right)Dq^2}{|\mathrm{\Omega }|\left(|\mathrm{\Omega }|+{\displaystyle \frac{2}{\tau _s}}\right)\left(|\mathrm{\Omega }|+Dq^2\right)}}\right]\left[\psi \left({\displaystyle \frac{1}{2}}+{\displaystyle \frac{1}{2\pi T_c\tau _s}}+{\displaystyle \frac{|\mathrm{\Omega }|}{2\pi T_c}}\right)\psi \left({\displaystyle \frac{1}{2}}+{\displaystyle \frac{1}{2\pi T_c\tau _s}}\right)\right]`$ (76)
$``$ $`{\displaystyle \frac{4(Dq^2)^2}{[\mathrm{\Omega }^2(Dq^2)^2]^2}}{\displaystyle \frac{\left[\psi \left({\displaystyle \frac{1}{2}}+{\displaystyle \frac{1}{2\pi T_c\tau _s}}+{\displaystyle \frac{|\mathrm{\Omega }|}{2\pi T_c}}\right)\psi \left({\displaystyle \frac{1}{2}}+{\displaystyle \frac{1}{2\pi T_c\tau _s}}\right)\right]^2}{\left[\psi \left({\displaystyle \frac{1}{2}}+{\displaystyle \frac{1}{2\pi T_c\tau _s}}+{\displaystyle \frac{Dq^2+|\mathrm{\Omega }|)}{4\pi T_c}}\right)\psi \left({\displaystyle \frac{1}{2}}+{\displaystyle \frac{1}{2\pi T_c\tau _s}}\right)\right]}}\}V_C(q,\mathrm{\Omega }).`$ (77)
Since the worst singularity possible in $`V_C(q,\mathrm{\Omega })`$ at $`q=0`$ goes as $`1/q^2`$, the overall $`q^2`$ factor multiplying $`V_C`$ in the above expression means that this singularity is removed. It follows that the removal of the $`q=0`$ singularity in the Coulomb part is unaffected by the addition of magnetic impurities – this is because it is a general feature enforced by gauge invariance, as we will show in Appendix B.
To calculate the Coulombic suppression term of Eqn. (75), we change variables to $`m=\mathrm{\Omega }/2\pi T`$ and $`y=Dq^2/2\pi T`$, noting that $`_q=𝑑y/(8\pi ^2DT)`$, and
$$N(0)V_C(q,\mathrm{\Omega })=N(0)\left[\frac{q^2}{4\pi e^2}+\frac{2N(0)Dq^2}{Dq^2+|\mathrm{\Omega }|}\right]^1\frac{Dq^2+|\mathrm{\Omega }|}{2Dq^2}=\frac{m+y}{2y}.$$
(78)
This leads to the result
$`\mathrm{log}\left({\displaystyle \frac{T_c}{T_{c0}}}\right)={\displaystyle \frac{1}{8\pi ^2N(0)D}}{\displaystyle \underset{m=1}{\overset{M}{}}}{\displaystyle _0^M}𝑑y`$ $`\{{\displaystyle \frac{1}{my}}\psi ^{}({\displaystyle \frac{1}{2}}+\alpha )`$ (79)
$`+`$ $`\left[{\displaystyle \frac{2y(m^2+y^2)}{m(m^2y^2)^2}}{\displaystyle \frac{2\alpha y}{m(m+2\alpha )(m+y)}}\right]\left[\psi \left({\displaystyle \frac{1}{2}}+\alpha +m\right)\psi \left({\displaystyle \frac{1}{2}}+\alpha \right)\right]`$ (80)
$``$ $`{\displaystyle \frac{4y}{(my)(m^2y^2)}}{\displaystyle \frac{\left[\psi \left(\frac{1}{2}+\alpha +m\right)\left(\frac{1}{2}+\alpha \right)\right]^2}{\left[\psi \left(\frac{1}{2}+\alpha +\frac{m+y}{2}\right)\psi \left(\frac{1}{2}+\alpha \right)\right]}}\},`$ (81)
where $`\alpha =1/2\pi T_c\tau _s`$ and the upper cutoff $`M=1/2\pi T_c\tau `$. The leading order term is that which goes like $`1/y`$ at large $`y`$, leading to logarithmic behavior. To isolate this, we add and subtract the term
$`{\displaystyle \underset{m=1}{\overset{M}{}}}{\displaystyle _0^M}{\displaystyle \frac{dy}{(m+y)}}\left\{\left({\displaystyle \frac{2}{m}}{\displaystyle \frac{2\alpha }{m(m+2\alpha )}}\right)\left[\psi \left({\displaystyle \frac{1}{2}}+\alpha +m\right)\psi \left({\displaystyle \frac{1}{2}}+\alpha \right)\right]\psi ^{}\left({\displaystyle \frac{1}{2}}+\alpha \right)\right\}`$ (82)
$`=`$ $`{\displaystyle \underset{m=1}{\overset{M}{}}}\mathrm{ln}\left({\displaystyle \frac{M+m}{m}}\right)\left\{{\displaystyle \frac{2(m+\alpha )}{m(m+2\alpha )}}\left[\psi \left({\displaystyle \frac{1}{2}}+\alpha +m\right)\psi \left({\displaystyle \frac{1}{2}}+\alpha \right)\right]\psi ^{}\left({\displaystyle \frac{1}{2}}+\alpha \right)\right\},`$ (83)
to give the result
$`\mathrm{ln}\left({\displaystyle \frac{T_c}{T_{c0}}}\right)=`$ $``$ $`{\displaystyle \frac{R_{\text{ }\text{ }\text{ }\text{ }}}{R_0}}\{{\displaystyle \underset{m=1}{\overset{M}{}}}\mathrm{ln}\left({\displaystyle \frac{M+m}{m}}\right)[{\displaystyle \frac{2(m+\alpha )}{m(m+2\alpha )}}[\psi ({\displaystyle \frac{1}{2}}+\alpha +m)\psi ({\displaystyle \frac{1}{2}}+\alpha )]\psi ^{}({\displaystyle \frac{1}{2}}+\alpha )]`$ (84)
$``$ $`{\displaystyle \underset{m=1}{\overset{M}{}}}{\displaystyle _0^M}dy{\displaystyle \frac{4y}{(my)(m^2y^2)}}[{\displaystyle \frac{\left[\psi \left(\frac{1}{2}+\alpha +m\right)\psi \left(\frac{1}{2}+\alpha \right)\right]^2}{\left[\psi \left(\frac{1}{2}+\alpha +\frac{m+y}{2}\right)\right]}}`$ (85)
$``$ $`[\psi ({\displaystyle \frac{1}{2}}+\alpha +m)\psi ({\displaystyle \frac{1}{2}}+\alpha )]+{\displaystyle \frac{ym}{2}}\psi ^{}({\displaystyle \frac{1}{2}}+\alpha )]\},`$ (86)
where we have noted that $`1/8\pi ^2N(0)D=R_{\text{ }\text{ }\text{ }\text{ }\text{ }}/R_0`$. We could now proceed to evaluate this expression, but before we do so, let us consider the domain of validity of the first-order perturbative result.
## V Beyond Perturbation Theory
Since we now have the full first-order perturbative correction to the transition temperature due to the effect of disorder on the Coulomb interaction, we could in principle plot curves of $`T_c(R_{\text{ }\text{ }\text{ }\text{ }\text{ }},1/\tau _s)`$ and compare to experiment. However the curves of $`T_c`$ vs $`R_{\text{ }\text{ }\text{ }\text{ }\text{ }}`$ for different values of $`1/\tau _s`$ would simply be exponential decays with different initial slopes. First-order perturbation theory is unable to treat the strong disorder region, and so cannot lead to the complete destruction of the superconductivity by non-magnetic disorder. If we are to consider the effects of arbitrary disorder strength, we must work beyond perturbation theory. In what follows we discuss two methods of doing this, and compare the results we obtain from them.
The simplest way to proceed is to “self-consistently” solve the first-order perturbative expression of Eq. (84). This simply means that we replace $`T_{c0}`$ by $`T_c`$ on the right-hand side of Eq. (84), and solve the implicit equation we obtain for $`T_c`$ which has the form
$$\mathrm{ln}\left(\frac{T_c}{T_{c0}}\right)=\psi \left(\frac{1}{2}\right)\psi \left(\frac{1}{2}+\frac{1}{2\pi T_c\tau _s}\right)\frac{R_{\text{ }\text{ }\text{ }\text{ }\text{ }}}{R_0}f(T_c,1/\tau _s).$$
(87)
Here $`f(T_c,1/\tau _s)`$ is the complicated expression on the right-hand side of Eq. (84), whilst the first term is just the mean-field suppression of $`T_c`$ by the magnetic impurities. From our knowledge of the situation without magnetic impurities, we know that unphysical re-entrance problems may arise with the solution of this equation, and we should not take it too seriously in the region where superconductivity is strongly suppressed.
The fact that we cannot trust the results obtained from this “self-consistent” theory leads us to ask the question of how to correctly go beyond first-order perturbation theory. The best approach is to derive the effective field theory from which the perturbation series may be deduced – in this case a non-linear sigma model – and treat this using the renormalization group. This has been done by Finkel’stein for the system without magnetic impurities, but has the problem that it is very difficult, and would become even more so if magnetic impurities were added. Recently Oreg and Finkel’stein demonstrated that the same results could be obtained using a much simpler non-perturbative resummation technique, which we show diagrammatically in Fig. (7). The method uses a featureless Coulomb interaction of magnitude $`N(0)V_C=1/2`$, consistent with the cancellation of the $`1/q^{d1}`$ divergence discussed earlier, and keeps only diagrams 3 and 4 of Fig. (6), since they give the greatest contribution. This leads to the equation for the pair scattering amplitude, $`\mathrm{\Gamma }(\omega _n,\omega _l)`$,
$$\mathrm{\Gamma }(\omega _n,\omega _l)=|\lambda |+t\mathrm{\Lambda }(\omega _n,\omega _l)\pi T\underset{m=(M+1)}{\overset{M}{}}[|\lambda |+t\mathrm{\Lambda }(\omega _n,\omega _m)]\frac{1}{|\omega _m|+1/\tau _s}\mathrm{\Gamma }(\omega _m,\omega _l),$$
(88)
where $`\omega _n=2\pi T(n+1/2)`$ is a Fermi Matsubara frequency, and the upper cut-off $`M=1/2\pi T\tau `$. The amplitude $`\mathrm{\Lambda }(\omega _n,\omega _l)`$ is given by
$$\mathrm{\Lambda }(\omega _n,\omega _l)=\{\begin{array}{cc}\mathrm{ln}\left[\frac{1}{(|\omega _n|+|\omega _l|)\tau }\right]\hfill & \omega _n\omega _l<0\hfill \\ \mathrm{ln}\left[\frac{1}{(|\omega _n|+|\omega _l|+2/\tau _s)\tau }\right]\hfill & \omega _n\omega _l>0\hfill \end{array}$$
(89)
where the breaking of time-reversal invariance by the spin-flip scattering means that $`\mathrm{\Lambda }`$ has a different form depending upon the relative signs of its two Matsubara frequencies. If we treat the $`\mathrm{\Gamma }(\omega _n,\omega _m)`$ as elements of a matrix $`\widehat{\mathrm{\Gamma }}`$, the matrix equation for $`\widehat{\mathrm{\Gamma }}`$ can be solved to yield
$$\widehat{\mathrm{\Gamma }}=\widehat{\omega }^{1/2}(\widehat{I}|\lambda |\widehat{\mathrm{\Pi }})^1\widehat{\omega }^{1/2}(|\lambda |\widehat{1}+t\widehat{\mathrm{\Lambda }}),$$
(90)
where
$$\widehat{\mathrm{\Pi }}=\frac{1}{2}\widehat{\omega }^{1/2}[\widehat{1}|\lambda |^1t\widehat{\mathrm{\Lambda }}]\widehat{\omega }^{1/2},$$
(91)
$`\widehat{\omega }_{nm}=(n+1/2+\alpha )\delta _{nm}`$, $`\widehat{\mathrm{\Lambda }}_{nm}=\mathrm{\Lambda }(\omega _n,\omega _m)`$, $`\widehat{1}_{nm}=1`$ and $`\widehat{I}_{nm}=\delta _{nm}`$. The matrix $`\widehat{\mathrm{\Gamma }}`$ becomes singular when an eigenvalue of $`\widehat{\mathrm{\Pi }}`$ equals $`1/|\lambda |`$, and this signals the onset of superconductivity. Note that the matrix $`\widehat{\mathrm{\Pi }}`$ depends on temperature both through the temperature dependence of its elements, and also through its rank $`2M`$. To find $`T_c`$, we start at the BCS value $`T_{c0}`$, which corresponds to a value of $`M`$ given by $`M_0=1/2\pi T_{c0}\tau `$. We decrease the temperature $`T`$ by increasing the upper cut-off $`M`$ successively by one. For each value of $`M`$, we construct the matrix $`\widehat{\mathrm{\Pi }}`$, and diagonalise it. When its lowest eigenvalue equals $`1/|\lambda |`$, we have reached the transition temperature $`T_c`$, which is given by $`T_c/T_{c0}=M_0/M`$. This method allows us to go to as low a temperature as we like, provided that we are prepared to diagonalize large enough matrices.
We will now plot curves of $`T_c`$ vs $`R_{\text{ }\text{ }\text{ }\text{ }\text{ }}`$ for fixed $`1/\tau _s`$, and $`T_c`$ vs $`1/\tau _s`$ for fixed $`R_{\text{ }\text{ }\text{ }\text{ }\text{ }}`$, derived both from the self-consistent perturbation theory of Eqn. (87), and from the non-perturbative resummation approach of Eqn. (88). This is done in Fig. (8), and we see that the two approaches are in rough agreement. The resummation technique is seen to remove the re-entrance problem which occurs in the $`\alpha =0`$ curve at large $`R_{\text{ }\text{ }\text{ }\text{ }\text{ }}`$, but surprisingly this re-entrance seems to be partially cured by the presence of magnetic impurities.
The above curves are fine from the theorist’s point of view, but experimentally what is measured is the suppression of $`T_c`$ by a certain fixed amount of magnetic impurities as the thickness of the superconductor is altered. The data is then presented in the form of the pair-breaking per magnetic impurity which can be written as
$$\alpha ^{}(R_{\text{ }\text{ }\text{ }\text{ }\text{ }})=\frac{T_c(R_{\text{ }\text{ }\text{ }\text{ }\text{ }},0)T_c(R_{\text{ }\text{ }\text{ }\text{ }\text{ }},1/\tau _s)}{1/\tau _s},$$
(92)
which we can also generate from our theoretical expressions. The result will of course depend upon the magnitude of the value of $`\alpha `$ we choose: we would like to choose $`\alpha `$ as small as possible so that we are always in the linear regime of pair-breaking, but not too small so that the difference is very sensitive to the discrete sums used in the numerical calculation. A typical plot is shown in Fig. (9). If we ignore the numerical noise we see that $`\alpha ^{}`$ is roughly constant and equal to its mean field value of $`\pi ^2/2`$. It only appears to increase as we approach the region where superconductivity is destroyed, and its total variation is only about $`10\%`$ even if we include this region. This is in agreement with the experimental data of Chervenak and Valles.
## VI Discussion and Conclusions
The main conclusion of this paper is that the effects of localization and interaction do not lead to an appreciable change of pair-breaking rate per magnetic impurity in disordered superconducting films provided that we are not too close to the superconductor-insulator transition. The experimental data agrees with this theoretical prediction, and thus confirms the validity of the basic model of $`T_c`$ suppression in disordered superconductors which consists of the BCS interaction, Coulomb repulsion and static disorder. The fact that the theoretical prediction is obtained both from first-order perturbation theory, and from a non-perturbative resummation technique, gives us increased confidence in its validity. Our calculations demonstrate that the resummation technique is a very powerful tool for going beyond perturbation theory which can be adapted to a variety of situations. Moreover we find that the ad hoc “self-consistent” extension of first-order perturbation theory can give sensible results even at values of $`R_{\text{ }\text{ }\text{ }\text{ }\text{ }}`$ near the superconductor-insulator transition, at least in the presence of pair-breaking.
The effect of nonmagnetic disorder on pair-breaking in superconducting films has previously been considered by Devereaux and Belitz using a model in which strong coupling effects are considered. Good agreement with experiment is also obtained with this approach, although more fitting parameters are required in this model. We note that only a single fitting parameter – the initial slope of the $`T_c(R_{\text{ }\text{ }\text{ }\text{ }\text{ }})`$ curve – is required in our approach. Unfortunately we see that the experimental data is unable to determine which, if any, of the two approaches is correct. In support of our approach we note that it is a “minimal model” in the sense that it contains the minimal physics to describe the system, and requires the input of a single fitting parameter. However, this is not to say that strong-coupling effects are not important in this system.
Another important result which emerges from the approach based on the grand-canonical potential is that the $`1/q^2`$ singularity of the disorder screened Coulomb potential is always cancelled in first-order perturbation theory. This removes the possibility of changing some experimental parameter to obtain a strong suppression of $`T_c`$ from this singularity. We have shown that this cancellation is enforced by gauge invariance, and leads us to suspect that it occurs to all orders in perturbation theory. It is this cancellation which makes it legitimate to use a featureless interaction in the resummation technique.
ACKNOWLEDGEMENTS
We thank I. Aleiner, A.M. Finkel’stein and Y. Oreg for helpful discussions. R.A.S. acknowledges support from the Nuffield Foundation. V.A. is supported by the U.S. National Science Foundation under grant DMR-9805613.
## A Calculation of Polarization Bubbles
In this appendix we give a detailed derivation of the polarization bubbles, $`\mathrm{\Pi }_{ij}`$, shown in Fig. (3). To evaluate these we must first calculate the impurity ladder, $`\mathrm{\Pi }`$, which is given by the geometric series
$$\mathrm{\Pi }=S+S\mathrm{\Gamma }_0S+S\mathrm{\Gamma }_0S\mathrm{\Gamma }_0S+\mathrm{},$$
(A1)
where $`\mathrm{\Gamma }_0`$ is the impurity line
$$\mathrm{\Gamma }_0=\frac{1}{2\pi N(0)\tau }\left[\lambda _1\tau _3\sigma _0\tau _3\sigma _0+\lambda _2\left(\tau _0\sigma _1\tau _0\sigma _1+\tau _0\sigma _2\tau _0\sigma _2+\tau _3\sigma _3\tau _3\sigma _3\right)\right],$$
(A2)
and $`1/\tau =1/\tau _0+1/\tau _s`$ is the total impurity scattering rate, $`\lambda _1=\tau /\tau _0`$, and $`\lambda _2=\tau /3\tau _s`$. $`S`$ is the momentum sum of a direct product of Green functions
$`S`$ $`=`$ $`{\displaystyle \underset{k}{}}G(k,i\omega )G(k+q,i\omega +i\mathrm{\Omega })`$ (A3)
$`=`$ $`\pi N(0)\tau I\left[\tau _3\sigma _0\tau _3\sigma _0{\displaystyle \frac{(\overline{z}\overline{\mathrm{\Delta }}\tau _1\sigma _3)(\overline{z}^{}\overline{\mathrm{\Delta }}^{}\tau _1\sigma _3)}{\epsilon \epsilon ^{}}}\right],`$ (A4)
and $`I`$ is the integral
$$I=\frac{1}{\pi \tau }𝑑\xi _k𝑑\widehat{\mathrm{\Omega }}\frac{\xi _k(\xi _k𝐪.𝐯_𝐅)}{(\xi _k^2\epsilon ^2)[(\xi _k𝐪.𝐯_𝐅)^2\epsilon ^2]}.$$
(A5)
Since we do not need the complete matrix structure of $`\mathrm{\Pi }`$, but just its traces with two matrices from the set $`\tau _1\sigma _3`$, $`\tau _2\sigma _3`$, $`\tau _3\sigma _0`$, we actually evaluate the impurity dressed vertices $`\mathrm{\Pi }_j`$ which have one matrix from the above set inserted between the two terms of the direct product in $`\mathrm{\Pi }`$. These satisfy the equation
$$\mathrm{\Pi }_j=S_j+S\mathrm{\Gamma }_0\mathrm{\Pi }_j.$$
(A6)
Starting with $`\mathrm{\Pi }_\mathrm{\Delta }`$ we see that
$`S_\mathrm{\Delta }`$ $`=`$ $`2\pi N(0)\tau I\left[\tau _3\sigma _0\tau _1\sigma _3\tau _3\sigma _0{\displaystyle \frac{(\overline{z}+\overline{\mathrm{\Delta }}\tau _1\sigma _3)\tau _1\sigma _3(\overline{z}^{}+\overline{\mathrm{\Delta }}^{}\tau _1\sigma _3)}{\epsilon \epsilon ^{}}}\right]`$ (A7)
$`=`$ $`2\pi N(0)\tau I\left[1{\displaystyle \frac{(\overline{z}+\overline{\mathrm{\Delta }}\tau _1\sigma _3)(\overline{z}^{}+\overline{\mathrm{\Delta }}^{}\tau _1\sigma _3)}{\epsilon \epsilon ^{}}}\right]\tau _1\sigma _3`$ (A8)
$`=`$ $`2\pi N(0)\tau I(\overline{\alpha }_+\beta _+\tau _1\sigma _3)\tau _1\sigma _3,`$ (A9)
where the $`\alpha `$ and $`\beta `$ terms are coherence factors
$$\alpha _\pm =1\frac{\overline{z}\overline{z}^{}\pm \overline{\mathrm{\Delta }}\overline{\mathrm{\Delta }}^{}}{\epsilon \epsilon ^{}};\overline{\alpha }_\pm =\alpha _\pm 2;\beta _\pm =\frac{\overline{z}^{}\overline{\mathrm{\Delta }}\pm \overline{z}\overline{\mathrm{\Delta }}^{}}{\epsilon \epsilon ^{}}.$$
(A10)
By inspection we see that $`\mathrm{\Pi }_\mathrm{\Delta }`$ must have the matrix form
$$\mathrm{\Pi }_\mathrm{\Delta }=2\pi N(0)\tau I[A+B\tau _1\sigma _3]\tau _1\sigma _3,$$
(A11)
and we now substitute this into Eqn. (A6) to deduce the coefficients $`A`$ and $`B`$. To derive the second term on the RHS of Eqn. (A6) we see that
$$\mathrm{\Gamma }_0\mathrm{\Pi }_\mathrm{\Delta }=I\lambda _1(AB\tau _1\sigma _3)\tau _1\sigma _3+3I\lambda _2(A+B\tau _1\sigma _3)\tau _1\sigma _3,$$
(A12)
and thus
$`S\mathrm{\Gamma }_0\mathrm{\Pi }_\mathrm{\Delta }`$ $`=`$ $`2\pi N(0)\tau I\{\lambda _1(A+B\tau _1\sigma _3)\tau _1\sigma _3+{\displaystyle \frac{(\overline{z}+\overline{\mathrm{\Delta }}\tau _1\sigma _3)\lambda _1(A+B\tau _1\sigma _3)\tau _1\sigma _3(\overline{z}^{}+\overline{\mathrm{\Delta }}^{}\tau _1\sigma _3)}{\epsilon \epsilon ^{}}}`$ (A14)
$`3\lambda _2(AB\tau _1\sigma _3)\tau _1\sigma _3{\displaystyle \frac{(\overline{z}+\overline{\mathrm{\Delta }}\tau _1\sigma _3)\lambda _2(AB\tau _1\sigma _3)\tau _1\sigma _3(\overline{z}^{}+\overline{\mathrm{\Delta }}^{}\tau _1\sigma _3)}{\epsilon \epsilon ^{}}}\}`$
$`=`$ $`2\pi N(0)\tau I\{(\lambda _13\lambda _2)A[1+{\displaystyle \frac{(\overline{z}+\overline{\mathrm{\Delta }}\tau _1\sigma _3)(\overline{z}^{}+\overline{\mathrm{\Delta }}^{}\tau _1\sigma _3)}{\epsilon \epsilon ^{}}}]\tau _1\sigma _3`$ (A16)
$`+(\lambda _1+3\lambda _2)B\tau _1\sigma _3[1{\displaystyle \frac{(\overline{z}+\overline{\mathrm{\Delta }}\tau _1\sigma _3)(\overline{z}^{}+\overline{\mathrm{\Delta }}^{}\tau _1\sigma _3)}{\epsilon \epsilon ^{}}}]\tau _1\sigma _3\}`$
$`=`$ $`2\pi N(0)\tau I[(\lambda _13\lambda _2)A(\overline{\alpha }_++\beta _+\tau _1\sigma _3)\tau _1\sigma _3+(\lambda _1+3\lambda _2)B(\alpha _+\beta _+\tau _1\sigma _3)].`$ (A17)
We can now equate the coefficients of $`1`$ and $`\tau _1\sigma _3`$ on the LHS and RHS of Eqn. (A6) to obtain the linear equations for $`A`$ and $`B`$,
$$\left[\begin{array}{cc}1+I(\lambda _13\lambda _2)\overline{\alpha }_+& I(\lambda _1+3\lambda _2)\beta _+\\ I(\lambda _13\lambda _2)\beta _+& 1I(\lambda _1+3\lambda _2)\alpha _+\end{array}\right]\left[\begin{array}{c}A\\ B\end{array}\right]=\left[\begin{array}{c}\overline{\alpha }_+\\ \beta _+\end{array}\right].$$
(A18)
This matrix equation can then be inverted by inverting the $`2\times 2`$ matrix and using the identity $`\alpha _+\overline{\alpha }_+=\beta _+^2`$ to obtain
$$\left[\begin{array}{c}A\\ B\end{array}\right]=\frac{1}{D_+}\left[\begin{array}{cc}1I(\lambda _1+3\lambda _2)\alpha _+& I(\lambda _1+3\lambda _2)\beta _+\\ I(\lambda _13\lambda _2)\beta _+& 1+I(\lambda _13\lambda _2)\overline{\alpha }_+\end{array}\right]\left[\begin{array}{c}\overline{\alpha }_+\\ \beta _+\end{array}\right]=\frac{1}{D_+}\left[\begin{array}{c}\overline{\alpha }_+\\ \beta _+\end{array}\right],$$
(A19)
where the determinant $`D_+`$ can be written
$`D_+`$ $`=`$ $`[1+I(\lambda _13\lambda _2)\overline{\alpha }_+][1I(\lambda _1+3\lambda _2)\alpha _+]+I^2(\lambda _1+3\lambda _2)(\lambda _13\lambda _2)\beta _+^2`$ (A20)
$`=`$ $`1+I(\lambda _13\lambda _2)\overline{\alpha }_+I(\lambda _1+3\lambda _2)\alpha _+=12I\lambda _1+6I\lambda _2\left({\displaystyle \frac{\overline{z}\overline{z}^{}+\overline{\mathrm{\Delta }}\overline{\mathrm{\Delta }}^{}}{\epsilon \epsilon ^{}}}\right).`$ (A21)
Similar results are obtained for $`\mathrm{\Pi }_\varphi `$ and $`\mathrm{\Pi }_\rho `$, leading to the results
$`\mathrm{\Pi }_\mathrm{\Delta }`$ $`=`$ $`2\pi N(0)\tau {\displaystyle \frac{I}{D_+}}(\overline{\alpha }_+\beta _+\tau _1\sigma _3)\tau _1\sigma _3`$ (A22)
$`\mathrm{\Pi }_\varphi `$ $`=`$ $`2\pi N(0)\tau {\displaystyle \frac{1}{D_{}}}(\overline{\alpha }_{}\beta _{}\tau _1\sigma _3)\tau _2\sigma _3`$ (A23)
$`\mathrm{\Pi }_\rho `$ $`=`$ $`2\pi N(0)\tau {\displaystyle \frac{I}{D_{}}}(\alpha _{}\beta _{}\tau _1\sigma _3)\tau _3\sigma _0,`$ (A24)
where
$$D_\pm =12I\lambda _1+6I\lambda _2\left(\frac{\overline{z}\overline{z}^{}\pm \overline{\mathrm{\Delta }}\overline{\mathrm{\Delta }}^{}}{\epsilon \epsilon ^{}}\right).$$
(A25)
If we evaluate the integral in Eqn. (A5) we find that $`I`$ is given by
$$2I\tau =\frac{1}{\overline{W}+\overline{W}^{}}\frac{q^2v_F^2}{2(\overline{W}+\overline{W}^{})^3},$$
(A26)
where
$$\overline{W}=\sqrt{\overline{\omega }^2+\overline{\mathrm{\Delta }}^2};\overline{W}=\sqrt{\omega ^2+\overline{\mathrm{\Delta }}^2}.$$
(A27)
From the second part of Eqn. (21) we see that we can write
$$\overline{W}+\overline{W}^{}=\frac{1}{\tau _0}\frac{1}{\tau _s}+\mathrm{\Delta }U+\mathrm{\Delta }U^{},$$
(A28)
and substituting Eqns. (A28) and (A26) into Eqn. (A25) gives the result
$$D_\pm =1\frac{\left[{\displaystyle \frac{1}{\tau _0}}{\displaystyle \frac{1}{\tau _s}}\left({\displaystyle \frac{uu^{}1}{UU^{}}}\right)\right]}{\left[{\displaystyle \frac{1}{\tau _0}}{\displaystyle \frac{1}{\tau _s}}+\mathrm{\Delta }U+\mathrm{\Delta }U^{}\right]}+Dq^2\tau \left[Dq^2+\mathrm{\Delta }U+\mathrm{\Delta }U^{}+\frac{1}{\tau _s}\left(\frac{uu^{}1}{UU^{}}1\right)\right]\tau .$$
(A29)
We can finally obtain the non-zero polarization bubbles $`\mathrm{\Pi }_{ij}`$ by inserting the second matrix from the set $`\tau _1\sigma _3`$, $`\tau _2\sigma _3`$, $`\tau _3\sigma _0`$ into $`\mathrm{\Pi }_j`$ and taking the trace. This yields
$`\mathrm{\Pi }_{\mathrm{\Delta }\mathrm{\Delta }}(q,\mathrm{\Omega })`$ $`=`$ $`\pi N(0)T{\displaystyle \underset{\omega }{}}\left[{\displaystyle \frac{UU^{}+uu^{}1}{UU^{}}}\right]{\displaystyle \frac{1}{\left(Dq^2+\mathrm{\Delta }U+\mathrm{\Delta }U^{}{\displaystyle \frac{1}{\tau _s}}\left[{\displaystyle \frac{UU^{}uu^{}+1}{UU^{}}}\right]\right)}}`$ (A30)
$`\mathrm{\Pi }_{\varphi \varphi }(q,\mathrm{\Omega })`$ $`=`$ $`\pi N(0)T{\displaystyle \underset{\omega }{}}\left[{\displaystyle \frac{UU^{}+uu^{}+1}{UU^{}}}\right]{\displaystyle \frac{1}{\left(Dq^2+\mathrm{\Delta }U+\mathrm{\Delta }U^{}{\displaystyle \frac{1}{\tau _s}}\left[{\displaystyle \frac{UU^{}uu^{}1}{UU^{}}}\right]\right)}}`$ (A31)
$`\mathrm{\Pi }_{\rho \rho }(q,\mathrm{\Omega })`$ $`=`$ $`\pi N(0)T{\displaystyle \underset{\omega }{}}\left[{\displaystyle \frac{UU^{}uu^{}1}{UU^{}}}\right]{\displaystyle \frac{1}{\left(Dq^2+\mathrm{\Delta }U+\mathrm{\Delta }U^{}{\displaystyle \frac{1}{\tau _s}}\left[{\displaystyle \frac{UU^{}uu^{}1}{UU^{}}}\right]\right)}}+N(0)`$ (A32)
$`\mathrm{\Pi }_{\varphi \rho }(q,\mathrm{\Omega })`$ $`=`$ $`\pi N(0)T{\displaystyle \underset{\omega }{}}\left[{\displaystyle \frac{u^{}u}{UU^{}}}\right]{\displaystyle \frac{1}{\left(Dq^2+\mathrm{\Delta }U+\mathrm{\Delta }U^{}{\displaystyle \frac{1}{\tau _s}}\left[{\displaystyle \frac{UU^{}uu^{}1}{UU^{}}}\right]\right)}}=\mathrm{\Pi }_{\rho \varphi }(q,\mathrm{\Omega }).`$ (A33)
## B Low-Momentum Singularities in Density and Phase Propagators
The identities
$$\lambda ^1+\mathrm{\Pi }_{\varphi \varphi }(0,\mathrm{\Omega })=x\mathrm{\Pi }_{\varphi \rho }(0,\mathrm{\Omega }),\mathrm{\Pi }_{\varphi \rho }(0,\mathrm{\Omega })=x\mathrm{\Pi }_{\rho \rho }(0,\mathrm{\Omega }),$$
(B1)
where $`x=\mathrm{\Omega }/2\mathrm{\Delta }`$, play a central role in the present paper and in I. As we have seen, their consequence
$$[\lambda ^1+\mathrm{\Pi }_{\varphi \varphi }(0,\mathrm{\Omega })]\mathrm{\Pi }_{\rho \rho }(0,\mathrm{\Omega })+\mathrm{\Pi }_{\varphi \rho }(0,\mathrm{\Omega })^2=0$$
(B2)
leads to the potentials $`V_{\varphi \varphi }`$, $`V_{\varphi \rho }`$, and $`V_{\rho \rho }`$ having $`1/q^{d1}`$ singularities at low momentum $`q`$ for all temperatures $`0TT_c`$ and all non-zero frequencies $`\mathrm{\Omega }0`$. The importance of these identities suggests that they embody an underlying invariance principle. In this appendix we show that they are Ward identities connected to charge conservation, which is very reasonable since the impossibility of instantaneously moving the conserved screening charge a finite distance is at the root of these singularites, The ideas at work here go back to Nambu’s 1960 paper and its elaborations and they are only included here for completeness. It is unfortunate that the physical basis for these identities was left obscure in I.
To avoid irrelevant notational complications, we shall work within the $`2\times 2`$ Nambu space. The $`4\times 4`$ space needed to deal with spin flip scattering does not affect the general argument, and we shall in any case explicitly verify the identities for this case later in this appendix.
The ‘proper polarization parts’, $`\mathrm{\Pi }`$, are calculated within a mean field approximation, in which the interactions are replaced according to
$$VV_{MF}=\mathrm{\Delta }Tr\mathrm{\Psi }^{}\tau _1\mathrm{\Psi },$$
(B3)
which implies a choice of phase for the order parameter. It is known that the quasiparticles obtained in this approximation do not conserve charge, because $`V_{MF}`$ does not commute with the electron density. Since the only other non-commuting part of the Hamiltonian is the kinetic energy, the operator equation of motion for the density $`\rho Tr\mathrm{\Psi }^{}\tau _3\mathrm{\Psi }`$ is
$$\frac{\rho }{t}+\stackrel{}{j}=i[V_{MF},\rho ]=2\mathrm{\Delta }Tr\mathrm{\Psi }^{}\tau _2\mathrm{\Psi },$$
(B4)
where $`\stackrel{}{j}=Tr[\mathrm{\Psi }^{}\stackrel{}{}\mathrm{\Psi }\mathrm{\Psi }\stackrel{}{}\mathrm{\Psi }^{}]`$ is the current density operator. Eq. (B4) leads to the identity
$`{\displaystyle \frac{}{t_2}}T\mathrm{\Psi }_i(\stackrel{}{x}_1,t_1)\rho (\stackrel{}{x}_2,t_2)\mathrm{\Psi }_j^{}(\stackrel{}{x}_3,t_3)=`$ (B5)
$`\delta (\stackrel{}{x}_1\stackrel{}{x}_2)\delta (t_1t_2)i[\tau _3G(\stackrel{}{x}_2,t_2,\stackrel{}{x}_3,t_3)]_{ij}\delta (\stackrel{}{x}_2\stackrel{}{x}_3)\delta (t_2t_3)i[G(\stackrel{}{x}_1,t_1,\stackrel{}{x}_2,t_2)\tau _3]_{ij}`$ (B6)
$`T\mathrm{\Psi }_i(\stackrel{}{x}_1,t_1)_2\stackrel{}{j}(\stackrel{}{x}_2,t_2)\mathrm{\Psi }_j^{}(\stackrel{}{x}_3,t_3)+T\mathrm{\Psi }_i(\stackrel{}{x}_1,t_1)\rho _\varphi (\stackrel{}{x}_2,t_2)\mathrm{\Psi }_j^{}(\stackrel{}{x}_3,t_3),`$ (B7)
where $`T`$ is the time ordering operator and we have defined $`\rho _\varphi Tr\mathrm{\Psi }^{}\tau _2\mathrm{\Psi }`$. \[The first two terms on the right of Eq. (B5) come from the derivative of the time ordering operator.\] Multiplying Eq. (B5) by $`\tau _3`$, taking the trace, and Fourier transforming in space and time leads in the limit of zero wave vector to the identity
$$\mathrm{\Omega }\mathrm{\Pi }_{\rho \rho }(0,\mathrm{\Omega })=2\mathrm{\Delta }\mathrm{\Pi }_{\rho \varphi }(0,\mathrm{\Omega }).$$
(B8)
On the other hand, multiplying by $`\tau _2`$ and performing these same operations yields
$`\mathrm{\Omega }\mathrm{\Pi }_{\varphi \rho }(0,\mathrm{\Omega })`$ $`=`$ $`{\displaystyle \frac{i}{\beta }}Tr[\tau _1G]+2\mathrm{\Delta }\mathrm{\Pi }_{\varphi \varphi }`$ (B9)
$`=`$ $`2\mathrm{\Delta }[\lambda ^1+\mathrm{\Pi }_{\varphi \varphi }].`$ (B10)
In the second line above the self consistency equation within the mean field approximation has been used. To obtain Eq. (B1) we must note that $`\mathrm{\Pi }_{\rho \varphi }`$ is antisymmetric in its indices because of time reversal invariance—under which $`\rho `$ is symmetric and $`\rho _\varphi `$ antisymmetric.
Since the impurity interaction commutes with the charge density these identities survive in any ‘conserving’ approximation and, in particular, in our sum of all non-overlapping graphs. In the main body of this paper, the mean field approximation was used as a way station on the road to the single loop approximation of Section III. There the phase of the order parameter is not fixed as in Eq. (B3) but determined self consistently, which restores charge conservation.
Finally, we shall verify explicitly that the identities (B1) are satisfied by our calculated expressions. We start with the equations for $`\lambda ^1+\mathrm{\Pi }_{\varphi \varphi }`$, $`\mathrm{\Pi }_{\varphi \rho }`$, and $`\mathrm{\Pi }_{\rho \rho }`$,
$`\lambda ^1+\mathrm{\Pi }_{\varphi \varphi }(0,\mathrm{\Omega })`$ $`=`$ $`T{\displaystyle \underset{\omega }{}}\left[{\displaystyle \frac{UU^{}+uu^{}+1}{UU^{}\left(U+U^{}\zeta \left[\frac{UU^{}uu^{}1}{UU^{}}\right]\right)}}{\displaystyle \frac{1}{U}}\right]`$ (B11)
$`\mathrm{\Pi }_{\varphi \rho }(0,\mathrm{\Omega })`$ $`=`$ $`T{\displaystyle \underset{\omega }{}}{\displaystyle \frac{u^{}u}{UU^{}\left(U+U^{}\zeta \left[\frac{UU^{}uu^{}1}{UU^{}}\right]\right)}}`$ (B12)
$`\mathrm{\Pi }_{\rho \rho }(0,\mathrm{\Omega })`$ $`=`$ $`{\displaystyle \frac{1}{\pi }}T{\displaystyle \underset{\omega }{}}{\displaystyle \frac{UU^{}uu^{}1}{UU^{}\left(U+U^{}\zeta \left[\frac{UU^{}uu^{}1}{UU^{}}\right]\right)}},`$ (B13)
where we have removed the common factor $`\pi N(0)`$ from each $`\mathrm{\Pi }`$, and set $`\mathrm{\Delta }=1`$ for algebraic convenience. These factors can, of course, be replaced when we have finished.
We will first prove the relationship between $`\lambda ^1+\mathrm{\Pi }_{\varphi \varphi }`$ and $`\mathrm{\Pi }_{\varphi \rho }`$, namely
$$\lambda ^1+\mathrm{\Pi }_{\varphi \varphi }(0,\mathrm{\Omega })=\frac{\mathrm{\Omega }}{2}\mathrm{\Phi }_{\varphi \rho }(0,\mathrm{\Omega }).$$
(B14)
To proceed note that we can write
$$uu^{}=\frac{1}{2}[u^2+u^2(u^{}u)^2]=\frac{1}{2}[U^2+U^22(u^{}u)^2],$$
(B15)
from which it follows that
$`UU^{}+uu^{}+1`$ $`=`$ $`{\displaystyle \frac{1}{2}}[(U+U^{})^2(u^{}u)^2]`$ (B16)
$`UU^{}uu^{}1`$ $`=`$ $`{\displaystyle \frac{1}{2}}[(U^{}U)^2(u^{}u)^2].`$ (B17)
In the last term on the RHS of Eq. (B11) for $`\mathrm{\Pi }_{\varphi \varphi }`$, we can use the transformation $`\omega \omega ^{}`$, under which the sum over $`\omega `$ is invariant. This leads to $`uu^{}`$ and $`UU^{}`$, so that
$$2T\underset{\omega }{}\frac{1}{U}=T\underset{\omega }{}\left[\frac{1}{U}+\frac{1}{U^{}}\right]=T\underset{\omega }{}\frac{U+U^{}}{UU^{}}.$$
(B18)
We can now write, using Eqs. (B11), (B16) and (B18),
$`2(\lambda ^1+\mathrm{\Pi }_{\varphi \varphi })\mathrm{\Omega }\mathrm{\Pi }_{\varphi \rho }`$ $`=`$ $`T{\displaystyle \underset{\omega }{}}\left[{\displaystyle \frac{(U+U^{})^2(u^{}u)^2+\mathrm{\Omega }(u^{}u)}{UU^{}\left(U+U^{}\zeta \left[\frac{UU^{}uu^{}1}{UU^{}}\right]\right)}}{\displaystyle \frac{U+U^{}}{UU^{}}}\right]`$ (B19)
$`=`$ $`T{\displaystyle \underset{\omega }{}}{\displaystyle \frac{(U+U^{})^2(u^{}u)(u^{}u\mathrm{\Omega })(U+U^{})\left(U+U^{}\zeta \left[\frac{UU^{}uu^{}1}{UU^{}}\right]\right)}{UU^{}\left(U+U^{}\zeta \left[\frac{UU^{}uu^{}1}{UU^{}}\right]\right)}}.`$ (B20)
From the definition of $`u`$ and $`u^{}`$ in Eqn. (22) we obtain the identity
$`u^{}u\mathrm{\Omega }`$ $`=`$ $`(u^{}\omega ^{})(u\omega )=\zeta \left({\displaystyle \frac{u^{}}{U^{}}}{\displaystyle \frac{u}{U}}\right)`$ (B21)
$`(u^{}u)(u^{}u\mathrm{\Omega })`$ $`=`$ $`\zeta (u^{}u)\left({\displaystyle \frac{u^{}}{U^{}}}{\displaystyle \frac{u}{U}}\right)`$ (B22)
$`=`$ $`\zeta \left[{\displaystyle \frac{u^2}{U^{}}}+{\displaystyle \frac{u^2}{U}}{\displaystyle \frac{uu^{}}{U^{}}}{\displaystyle \frac{uu^{}}{U}}\right]`$ (B23)
$`=`$ $`\zeta \left[U+U^{}{\displaystyle \frac{1}{U}}{\displaystyle \frac{1}{U^{}}}{\displaystyle \frac{uu^{}}{U}}{\displaystyle \frac{uu^{}}{U^{}}}\right]`$ (B24)
$`=`$ $`\zeta (U+U^{})\left[{\displaystyle \frac{UU^{}uu^{}1}{UU^{}}}\right].`$ (B25)
It follows that the numerator in Eqn. (B19) is zero, and hence we have proved the required result (B14).
We next prove the relation between $`\mathrm{\Pi }_{\rho \rho }`$ and $`\mathrm{\Pi }_{\varphi \rho }`$, namely
$$\mathrm{\Pi }_{\varphi \rho }(0,\mathrm{\Omega })=\frac{\mathrm{\Omega }}{2}\mathrm{\Pi }_{\rho \rho }(0,\mathrm{\Omega }).$$
(B26)
We start by considering the sum
$$\underset{\omega _0\mathrm{}}{lim}T\underset{\omega =(\omega _0+\mathrm{\Omega })}{\overset{\omega _0}{}}\frac{u}{\mathrm{\Omega }U}.$$
(B27)
As $`|\omega |\mathrm{}`$, $`u(|\omega |+\zeta )\text{sgn}(\omega )`$ and $`u/U\text{sgn}(\omega )`$. Since there are $`\mathrm{\Omega }/2\pi T`$ more negative terms than positive, the sum becomes
$$\underset{\omega _0\mathrm{}}{lim}T\underset{\omega =(\omega _0+\mathrm{\Omega })}{\overset{\omega _0}{}}\frac{u}{\mathrm{\Omega }U}=\left(\frac{\mathrm{\Omega }}{2\pi T}\right)\left(\frac{T}{\mathrm{\Omega }}\right)=\frac{1}{2\pi }.$$
(B28)
We have chosen the limits so that we are able to make the usual $`\omega \omega ^{}`$ transformation in this sum. It follows that
$$\frac{1}{\pi }=2T\underset{\omega }{}\frac{u}{U\mathrm{\Omega }}=T\underset{\omega }{}\frac{1}{\mathrm{\Omega }}\left(\frac{u^{}}{U^{}}\frac{u}{U}\right)=T\underset{\omega }{}\frac{U+U^{}}{\mathrm{\Omega }(u^{}u)}\left(\frac{UU^{}uu^{}1}{UU^{}}\right),$$
(B29)
where we first make the $`\omega \omega ^{}`$ transformation, and then use (B21). We can then rewrite Eq. (B11) for $`\mathrm{\Pi }_{\rho \rho }`$ in the form
$`\mathrm{\Pi }_{\rho \rho }`$ $`=`$ $`T{\displaystyle \underset{\omega }{}}\left({\displaystyle \frac{UU^{}uu^{}1}{UU^{}}}\right)\left[{\displaystyle \frac{U+U^{}}{\mathrm{\Omega }(u^{}u)}}{\displaystyle \frac{1}{\left(U+U^{}\zeta \left[\frac{UU^{}uu^{}1}{UU^{}}\right]\right)}}\right]`$ (B30)
$`=`$ $`{\displaystyle \frac{T}{\mathrm{\Omega }}}{\displaystyle \underset{\omega }{}}\left({\displaystyle \frac{UU^{}uu^{}1}{UU^{}}}\right){\displaystyle \frac{(U+U^{})\left(U+U^{}\zeta \left[\frac{UU^{}uu^{}1}{UU^{}}\right]\right)(uu^{})\mathrm{\Omega }}{(u^{}u)\left(U+U^{}\zeta \left[\frac{UU^{}uu^{}1}{UU^{}}\right]\right)}}.`$ (B31)
From the identity (B21) we see that the numerator of (B30) can be rewritten as
$`(U+U^{})\left(U+U^{}\zeta \left[{\displaystyle \frac{UU^{}uu^{}1}{UU^{}}}\right]\right)(u^{}u)\mathrm{\Omega }`$ (B32)
$`=`$ $`(U+U^{})^2(u^{}u)(u^{}u\mathrm{\Omega })(u^{}u)\mathrm{\Omega }`$ (B33)
$`=`$ $`(U+U^{})^2(u^{}u)^2,`$ (B34)
and inserting the identity (B21) into the first factor in (B30), we get
$$\mathrm{\Pi }_{\rho \rho }=\frac{T}{2\mathrm{\Omega }}\underset{\omega }{}\frac{[(U^{}U)^2(u^{}u)^2][(U+U^{})^2(u^{}u)^2]}{(u^{}u)UU^{}\left(U+U^{}\zeta \left[\frac{UU^{}uu^{}1}{UU^{}}\right]\right)}.$$
(B35)
Multiplying out the numerator yields
$`[(U^{}U)^2(u^{}u)^2][(U+U^{})^2(u^{}u)^2]`$ (B36)
$`=`$ $`(u^2u^2)^22(u^{}u)^2(U^2+U^2)+(u^{}u)^4`$ (B37)
$`=`$ $`(u^{}u)^2[(u+u^{})^22u^22u^24+(u^{}u)^2]`$ (B38)
$`=`$ $`4(u^{}u)^2,`$ (B39)
from which it follows that
$$\mathrm{\Pi }_{\rho \rho }(0,\mathrm{\Omega })=\frac{2T}{\mathrm{\Omega }}\underset{\omega }{}\frac{u^{}u}{UU^{}\left(U+U^{}\zeta \left[\frac{UU^{}uu^{}1}{UU^{}}\right]\right)}=\frac{2}{\mathrm{\Omega }}\mathrm{\Pi }_{\varphi \rho }(0,\mathrm{\Omega }),$$
(B40)
completing our proof of the result (B26).
## C Evaluating Derivatives of $`\mathrm{\Pi }_{ij}`$ with respect to $`\mathrm{\Delta }`$
In this appendix we evaluate the derivatives of the polarization bubbles $`\mathrm{\Pi }_{ij}`$ with respect to the order parameter $`\mathrm{\Delta }`$ so that we may evaluate the first order correction to the order parameter self-consistency equation. The formulas for the $`\mathrm{\Pi }_{ij}`$ are given in Eqn. (A30), and we see that the derivative can operate either on the coherence factor or the denominator present in these expressions.
The difficulty in evaluating these derivatives arises because $`u(\omega )`$ satisfies the transcendental equation
$$\frac{\omega }{\mathrm{\Delta }}=u\left[1\frac{1}{\mathrm{\Delta }\tau _s}\frac{1}{(u^2+1)^{1/2}}\right],$$
(C1)
from which it follows that
$$\frac{\omega }{\mathrm{\Delta }^2}=\frac{u}{\mathrm{\Delta }}\left[1\frac{1}{\mathrm{\Delta }\tau _s}\frac{1}{(u^2+1)^{3/2}}\right]+\frac{1}{\mathrm{\Delta }^2\tau _s}\frac{u}{(u^2+1)^{1/2}},$$
(C2)
and thus
$$\frac{u}{\mathrm{\Delta }}=\frac{u}{\mathrm{\Delta }}\left[1\frac{1}{\mathrm{\Delta }\tau _s}\frac{1}{(u^2+1)^{3/2}}\right]^1,$$
(C3)
with a similar result for $`u^{}`$.
We first consider the effect of $`/\mathrm{\Delta }`$ on the coherence factors present in the $`\mathrm{\Pi }_{ij}`$. We see that it suffices to evaluate the derivatives
$$\frac{}{\mathrm{\Delta }}\{\frac{uu^{}}{UU^{}};\frac{1}{UU^{}};\frac{u^{}u}{UU^{}}\}.$$
(C4)
The first term in Eqn. (C4) gives
$$\frac{}{\mathrm{\Delta }}\left[\frac{uu^{}}{UU^{}}\right]=\frac{u^{}}{U^{}}[\frac{1}{U}\frac{u^2}{U^3}]\frac{u}{\mathrm{\Delta }}+(uu^{})=\frac{1}{\mathrm{\Delta }}\frac{uu^{}}{U^3U^{}}[1\frac{\zeta }{U^3}]^1+(uu^{}).$$
(C5)
Since this expression will occur inside a sum over $`\omega `$, and will multiply an expression that is invariant under the transformation $`\omega \omega ^{}`$, we see that the two terms in Eqn. (C5) will give equal results. Thus
$$\frac{}{\mathrm{\Delta }}\left[\frac{uu^{}}{UU^{}}\right]\frac{2}{\mathrm{\Delta }}\frac{uu^{}}{U^3U^{}}\left[1\frac{\zeta }{U^3}\right]^1.$$
(C6)
The second term in Eqn. (C4) gives
$$\frac{}{\mathrm{\Delta }}\left[\frac{1}{UU^{}}\right]=\frac{1}{U^{}}\frac{u}{U^3}\frac{u}{\mathrm{\Delta }}+(uu^{})\frac{2}{\mathrm{\Delta }}\frac{u^2}{UU^{}}[1\frac{\zeta }{U^3}]^1,$$
(C7)
whilst the third term in Eqn. (C4) gives
$$\frac{}{\mathrm{\Delta }}\left[\frac{u^{}u}{UU^{}}\right]=\frac{u^{}}{U^{}}\frac{u}{U^3}\frac{u}{\mathrm{\Delta }}\frac{1}{U^{}}\frac{1}{U^3}\frac{u}{\mathrm{\Delta }}+(uu^{})\frac{2}{\mathrm{\Delta }}\frac{u(uu^{}+1)}{U^3U^{}}[1\frac{\zeta }{U^3}]^1.$$
(C8)
The effect of $`/\mathrm{\Delta }`$ on the coherence factors can then be summarised in the form
$`{\displaystyle \frac{}{\mathrm{\Delta }}}\left[{\displaystyle \frac{uu^{}\pm 1}{UU^{}}}\right]`$ $``$ $`{\displaystyle \frac{2}{\mathrm{\Delta }}}{\displaystyle \frac{u(u^{}u)}{U^3U^{}}}\left[1{\displaystyle \frac{\zeta }{U^3}}\right]^1`$ (C9)
$`{\displaystyle \frac{}{\mathrm{\Delta }}}\left[{\displaystyle \frac{u^{}u}{UU^{}}}\right]`$ $``$ $`{\displaystyle \frac{2}{\mathrm{\Delta }}}{\displaystyle \frac{u(uu^{}+1)}{U^3U^{}}}\left[1{\displaystyle \frac{\zeta }{U^3}}\right]^1.`$ (C10)
Next we must consider the effect of $`/\mathrm{\Delta }`$ on the two denominators
$$D_\pm =Dq^2+\mathrm{\Delta }U+\mathrm{\Delta }U^{}\frac{1}{\tau _s}\left(1\frac{uu^{}1}{UU^{}}\right).$$
(C11)
The last term on the RHS is a coherence factor, and its derivative can be read off from Eqn. (C9) above. The only terms then to consider are $`\mathrm{\Delta }U`$ and $`\mathrm{\Delta }U^{}`$, which, of course, will give identical results after summation over $`\omega `$. We see that
$$\frac{}{\mathrm{\Delta }}(\mathrm{\Delta }U)=U+\mathrm{\Delta }\frac{u}{U^3}\frac{u}{\mathrm{\Delta }}=U\frac{u^2}{U^3}\left[1\frac{\zeta }{U^3}\right]^1=\frac{1}{U}\left[1\frac{\zeta }{U}\right]\left[1\frac{\zeta }{U^3}\right]^1.$$
(C12)
From this we obtain the final result
$$\frac{}{\mathrm{\Delta }}D_\pm =\left\{\frac{1}{U}\frac{\zeta }{U^2}\left(1+\frac{u(u^{}u)}{UU^{}}\right)\right\}\left[1\frac{\zeta }{U^3}\right]^1.$$
(C13)
Having now evaluated the action of $`/\mathrm{\Delta }`$ on all the components of the polarization bubbles, $`\mathrm{\Pi }_{ij}`$, we can now write down the results for the $`\mathrm{\Pi }_{ij}/\mathrm{\Delta }`$,
$`{\displaystyle \frac{\mathrm{\Pi }_{\mathrm{\Delta }\mathrm{\Delta }}}{\mathrm{\Delta }}}`$ $`=`$ $`2\pi N(0)T{\displaystyle \underset{\omega }{}}\left\{\left[1{\displaystyle \frac{\zeta }{U^3}}\right]^1\times {\displaystyle \frac{1}{\mathrm{\Delta }}}{\displaystyle \frac{u(u^{}+u)}{U^3U^{}}}{\displaystyle \frac{1}{D_+}}+\left(1+{\displaystyle \frac{uu^{}1}{UU^{}}}\right)\left\{{\displaystyle \frac{1}{U}}{\displaystyle \frac{\zeta }{U^2}}\left(1+{\displaystyle \frac{u(u^{}u)}{UU^{}}}\right)\right\}{\displaystyle \frac{1}{D_+^2}}\right\}`$ (C14)
$`{\displaystyle \frac{\mathrm{\Pi }_{\varphi \varphi }}{\mathrm{\Delta }}}`$ $`=`$ $`2\pi N(0)T{\displaystyle \underset{\omega }{}}\left\{\left[1{\displaystyle \frac{\zeta }{U^3}}\right]^1\times {\displaystyle \frac{1}{\mathrm{\Delta }}}{\displaystyle \frac{u(u^{}u)}{U^3U^{}}}{\displaystyle \frac{1}{D_{}}}+\left(1+{\displaystyle \frac{uu^{}+1}{UU^{}}}\right)\left\{{\displaystyle \frac{1}{U}}{\displaystyle \frac{\zeta }{U^2}}\left(1+{\displaystyle \frac{u(u^{}+u)}{UU^{}}}\right)\right\}{\displaystyle \frac{1}{D_{}^2}}\right\}`$ (C15)
$`{\displaystyle \frac{\mathrm{\Pi }_{\rho \rho }}{\mathrm{\Delta }}}`$ $`=`$ $`2\pi N(0)T{\displaystyle \underset{\omega }{}}\left\{\left[1{\displaystyle \frac{\zeta }{U^3}}\right]^1\times {\displaystyle \frac{1}{\mathrm{\Delta }}}{\displaystyle \frac{u(u^{}u)}{U^3U^{}}}{\displaystyle \frac{1}{D_{}}}\left(1{\displaystyle \frac{uu^{}+1}{UU^{}}}\right)\left\{{\displaystyle \frac{1}{U}}{\displaystyle \frac{\zeta }{U^2}}\left(1+{\displaystyle \frac{u(u^{}+u)}{UU^{}}}\right)\right\}{\displaystyle \frac{1}{D_{}^2}}\right\}`$ (C16)
$`{\displaystyle \frac{\mathrm{\Pi }_{\varphi \rho }}{\mathrm{\Delta }}}`$ $`=`$ $`2\pi N(0)T{\displaystyle \underset{\omega }{}}\left\{\left[1{\displaystyle \frac{\zeta }{U^3}}\right]^1\times {\displaystyle \frac{1}{\mathrm{\Delta }}}{\displaystyle \frac{u(uu^{}+1)}{U^3U^{}}}{\displaystyle \frac{1}{D_{}}}\left({\displaystyle \frac{u^{}u}{UU^{}}}\right)\left\{{\displaystyle \frac{1}{U}}{\displaystyle \frac{\zeta }{U^2}}\left(1+{\displaystyle \frac{u(u^{}+u)}{UU^{}}}\right)\right\}{\displaystyle \frac{1}{D_{}^2}}\right\}.`$ (C17)
FIGURES |
warning/0002/gr-qc0002075.html | ar5iv | text | # On the experimental foundations of the Maxwell equations
## 1 Introduction
According to special relativity and to metric theories of gravity like general relativity, the dynamics of the electromagnetic field is intimately connected with the structure of spacetime. For example, light cones can be used to identify the causal structure of spacetime in these theories, and the dynamics of the electromagnetic field combines with the quantum mechanics of charged particles to determine the behavior of atomic clocks and standards of length one uses to map out spacetime geometry. This connection between spacetime structure and the dynamics of the electromagnetic field clearly motivates the sharpest possible experimental tests of the validity of the familiar Maxwellian dynamics.
Generalized Maxwell equations provide a context in which to design and interpret such tests. Section 2 reviews the derivation of these equations from an operational definition of the electromagnetic field and the most basic notions of what constitutes a dynamical field theory. Since no geometrical spacetime structure like metric or conmnection is presumed by the derivation the generalized Maxwell equations can predict birefringence, charge non-conservation, wave damping and other effects not predicted by the familiar Maxwell equations. The familiar equations are, however, a special case of the generalized ones so that experiments which search for effects like those just mentioned can determine the extent to which that the dynamics of the electromagnetic field is or is not compatible with spacetime geometry. Section 3 discusses astronomical observations whose results constrain electromagnetic field dynamics that predict finite wave propagation speeds to be very close to the familiar Maxwellian dynamics. It also discusses effects which could be used to test the assumption of finite wave propagation speed and the assumption that the electrodynamic field solves a well-posed Cauchy problem. Section 4 contains a few closing remarks.
## 2 A derivation of the Maxwell equations
The mathematical structure of the Maxwell equations have been analyzed in and . A way to derive the Maxwell equations by means of an operational approach introduces the electromagnetic field $`F`$ via consideration of charged particle interferometry. Since the phase shift for small areas is found to be proportional to an interferometer’s area, $`\varphi `$ area. The proportionality factor can be defined to be the electromagnetic field $`F`$: $`\varphi =\frac{1}{2}F_{\mu \nu }\sigma ^{\mu \nu }`$ (the area is described by an antisymmetric second rank tensor). In this way we operationally defined the electromagnetic field.
The experimental uniqueness of the phase shift immediately requires
$$0=dF.$$
(1)
In a $`3+1`$–slicing of the manifold, these equations consist in three dynamical equations $`\mathrm{\pounds }_v\overline{B}=\overline{d}\overline{E}`$ ($`v`$ is a “timelike” vector field) $`\overline{d}\overline{B}=0`$ . Another three dynamical equations are needed.
These other equations can be obtained by requiring that the Maxwell field obey an evolutional structure with respect to the slicing. This is equivalent to requiring an evolution equation $`\mathrm{\pounds }_vF=\lambda (F)`$ that is compatible with (1). In coordinates: $`\mathrm{\pounds }_vF_{\widehat{\mu }\widehat{\nu }}=\lambda _{\widehat{\mu }\widehat{\nu }}(F)`$, $`\mathrm{\pounds }_vF_{\mu 0}=\lambda _\mu (F)=\lambda _\mu ^0(F)+j_\mu `$, where spatial indices are denoted by a hat. The first equation is contained in (1). We call the part of $`\lambda _\mu (F)`$ which remains as $`F0`$ the source $`j`$ of electromagnetic field. The requirement that the superposition principle should hold leads to a linear evolution equation as well as to constraints which are fulfilled by the superposition, too.
The additional requirement that the electromagnetic field should propagate with a finite propagation speed implies that $`\lambda `$ is a differential operator and that it should be of first order only . That gives our generalized Maxwell equations the form
$$_{[\rho }F_{\mu \nu ]}=0_\nu (\lambda _\mu ^{\nu \rho \sigma }F_{\rho \sigma })+\overline{\lambda }_\mu ^{\rho \sigma }F_{\rho \sigma }=j_\mu .$$
(2)
In general, these equations show birefringence, that is, two light cones, which of course violate Local Lorentz Invariance (LLI). The requirement of a unique light cone implies the existence of a non–degenerate second rank tensor $`g^{\mu \nu }`$ so that $`\lambda ^{\mu \nu \rho \sigma }`$ can be rewritten as
$$\lambda ^{\mu \nu \rho \sigma }=\frac{1}{2}\sqrt{g}(g^{\rho [\mu }g^{\nu ]\sigma }g^{\sigma [\mu }g^{\nu ]\rho })+\theta ϵ^{\mu \nu \rho \sigma }$$
(3)
where $`ϵ^{\mu \nu \rho \sigma }`$ is the total antisymmetric symbol. Therefore we can rewrite the inhomogeneous Maxwell equations (2) as ($`D`$ is the Riemannian covariant derivative)
$$j^\nu =D_\mu F^{\mu \nu }+\mathrm{\Sigma }_{\mu \rho }^\nu F^{\mu \rho }\text{with}\mathrm{\Sigma }_{\mu \rho }^\nu :=\overline{\lambda }_{\mu \rho }^\nu ϵ^{\nu \sigma }{}_{\mu \rho }{}^{}_{\sigma }^{}\theta \{_{\sigma [\mu }^\sigma \}\delta _{\rho ]}^\nu .$$
(4)
The form of $`\mathrm{\Sigma }`$ can be restricted by requiring charge conservation in the form $`D_\mu j^\mu =0`$. The most general $`\mathrm{\Sigma }`$ compatible with this requirement is
$$\mathrm{\Sigma }_{\mu \nu \rho }=ϵ_{\mu \nu \rho }{}_{}{}^{\sigma }_{\sigma }^{}\theta .$$
(5)
Ni has shown that this coupling does not violate the weak equivalence principle for falling charges. In addition, it was established in that vacuum polarization effects of the Dirac equation in a space–time with torsion leads to an effective Maxwell equation with an additional coupling of the above form. $`\theta `$ can be interpreted as axion field which appears in the low–energy limit of string theory and is a candidate for the dark matter in the universe. These effective Maxwell equations can be obtained from a Lagrangian .
The coupling to $`\theta `$ results in a violation of Local Position Invariance. However, in a rotating and accelerating frame with rotation $`𝝎`$ and acceleration $`𝒂`$ (stationary situation, Newtonian potential $`U(0)=0`$) the above Maxwell equations acquire the form (with the PPN parameter $`\gamma `$)
$`j^0`$ $`=`$ $`\mathbf{}𝑬(𝒂(1+\gamma )\mathbf{}U)𝑬2(𝝎\mathbf{}\theta )𝑩`$
$`𝒋`$ $`=`$ $`\mathbf{}\times 𝑩(𝒂(\gamma +1)\mathbf{}U)\times 𝑩+2(𝝎\mathbf{}\theta )\times 𝑬`$ (6)
which means that gravity and the pseudoscalar field can be transformed away. The effects of $`\mathbf{}\theta `$ or axial torsion can be simulated by a rotation of the reference frame .
The coupling to $`\mathbf{}\theta `$ amounts to a non–metric Faraday–effect: If at the worldline of the source the polarization of the electromagnetic field is parallelly propagated, $`D_uF=0`$, then the observer sees a rotation of the polarization along his worldline, $`D_vF_{\mu \nu }=2\dot{x}k^\rho (_\tau \theta )ϵ_{\rho [\mu }^{\sigma \tau }F_{\nu ]\sigma }`$. The absence of such an effect implies usual Maxwell equations coupled to a space–time metric only:
$$dF=0dF=j.$$
(7)
## 3 Experiments testing the Maxwell equations
In principle, observation of any electromagnetic phenomenon can provide a basis for a test of the validity of the Maxwell equations, but phenomena associated with wave propagation are particularly well suited to imposing sharp constraints on the dynamics of the electromagnetic field since astronomical observations can sometimes impose limits on effects that have built up in the course of propagation over huge distances. Laboratory tests can be considered when circumstances do not allow astronomical observations to distinguish between effects caused by departures from familiar Maxwell dynamics and by other physical processes.
### 3.1 Tests of the generalized Maxwell equations
First we derive a wave equation from the generalized Maxwell equations (2)
$$_{[\mu }j_{\nu ]}=\delta _{[\mu }^\alpha \lambda _{\nu ]}{}_{}{}^{\kappa \rho \sigma }_{\alpha }^{}_\kappa F_{\rho \sigma }+\stackrel{~}{\lambda }_{[\nu \mu ]}{}_{}{}^{\kappa \rho \sigma }_{\kappa }^{}F_{\rho \sigma }+\widehat{\lambda }_{[\nu \mu ]}{}_{}{}^{\rho \sigma }F_{\rho \sigma }^{},$$
(8)
where $`\stackrel{~}{\lambda }`$ and $`\widehat{\lambda }`$ are combinations of the original $`\lambda `$ and $`\overline{\lambda }`$ and its derivatives. If $`\lambda _\mu {}_{}{}^{\nu \rho \sigma }=\frac{1}{2}(\delta _\mu ^\rho g^{\nu \sigma }\delta _\mu ^\nu g^{\rho \sigma })+\delta \lambda _\mu ^{\nu \rho \sigma }`$ and if the derivatives of the coefficients are assumed to be negligible, then we get from (8) for $`j=0`$
$$0=\left(\delta _\mu ^\rho \delta _\nu ^\sigma \mathrm{}+2\delta \lambda _{[\mu }{}_{}{}^{\kappa \rho \sigma }_{\nu ]}^{}_\kappa \right)F_{\rho \sigma }+2\overline{\lambda }_{[\mu }{}_{}{}^{\rho \sigma }_{\nu ]}^{}F_{\rho \sigma }.$$
(9)
We assume that $`\delta \lambda _\nu ^{\nu \rho \sigma }`$ as well as $`\overline{\lambda }_\nu ^{\rho \sigma }`$ are small.
#### 3.1.1 Propagation of characteristics
First we discuss the behavior of characteristics, or shock waves, predicted by (9). This is determined by this equation’s principal polynomial ($`k`$ is the normal of the characteristic surface)
$$0=det\left(\delta _{[\mu }^\rho \delta _{\nu ]}^\sigma k^2+2\delta \lambda _{[\mu }{}_{}{}^{\kappa \rho \sigma }k_{\nu ]}^{}k_\kappa \right).$$
(10)
This prediction and the results of astrophysical observations that search for birefringence effects impose the constraint $`\delta \lambda _\mu {}_{}{}^{\kappa \rho \sigma }10^{28}`$ .
#### 3.1.2 Propagation of plane waves
If we take the electromagnetic field to be a plane wave $`F=F_0e^{ikx}`$, and do not take the eikonal limit, then the equation (9) implies
$$0=det\left(\delta _{[\mu }^\rho \delta _{\nu ]}^\sigma k^2+2\delta \lambda _{[\mu }{}_{}{}^{\kappa \rho \sigma }k_{\nu ]}^{}k_\kappa +2i\overline{\lambda }_{[\mu }{}_{}{}^{\rho \sigma }k_{\nu ]}^{}\right).$$
(11)
The factor $`i`$ associated with the $`\overline{\lambda }`$ stems from a term involving a single field derivative. To first order in $`\delta \lambda `$ and $`\overline{\lambda }`$ the corresponding dispersion relation is ($`det(1+\mathrm{\Lambda })=1+\text{tr}\mathrm{\Lambda }+𝒪(\mathrm{\Lambda }^2)`$)
$$0=(k^2)^5\left(k^2+\delta \lambda _{[\mu }{}_{}{}^{\kappa \mu \nu }k_{\nu ]}^{}k_\kappa +\overline{\lambda }_{[\mu }{}_{}{}^{\mu \nu }k_{\nu ]}^{}\right)+𝒪(\delta \lambda ^2,\overline{\lambda }^2).$$
(12)
Non–trivial solutions are given by ($`k^2=\omega ^2\mathrm{k}^2`$)
$`\omega _{1,2}`$ $`=`$ $`{\displaystyle \frac{1}{2}}((\delta \lambda _\mu {}_{}{}^{0\mu \widehat{\nu }}+\delta \lambda _\mu {}_{}{}^{\widehat{\kappa }\nu 0})k_{\widehat{\nu }}+i\overline{\lambda }_\mu {}_{}{}^{\mu 0})`$ (13)
$`\pm (\mathrm{k}(1{\displaystyle \frac{1}{2}}\delta \lambda _\mu {}_{}{}^{0\mu 0}){\displaystyle \frac{1}{2}}(\delta \lambda _\rho {}_{}{}^{\widehat{\mu }\rho \widehat{\nu }}{\displaystyle \frac{k_{\widehat{\mu }}}{\mathrm{k}}}k_{\widehat{\nu }}+i\overline{\lambda }_\mu {}_{}{}^{\mu \widehat{\nu }}{\displaystyle \frac{k_{\widehat{\nu }}}{\mathrm{k}}})).`$
Note the presence of two effects: (i) an anisotropic relation between $`k_{\widehat{\mu }}`$ and $`\omega `$ due to $`\delta \lambda `$, indicating a violation of LLI, and (ii) imaginary terms due to a trace of $`\overline{\lambda }`$, indicating a damping of a plane wave’s intensity. This damping is independent of wave length but depends on the direction in which a wave propagates. Knowledge of the distance and intrinsic brightness of stars should permit the estimation of the strength of any such damping. In principle, tests of LLI can lead to estimates of $`\delta \lambda `$ and observations of damping lead to estimates of $`\overline{\lambda }`$ and, thus, constraints on charge non-conservation. Recall that charge conservation forced $`\overline{\lambda }=0`$.
In the case of a coupling to the axion alone we find the exact dispersion relation $`\omega =\pm \mathrm{k}\sqrt{1\pm {\displaystyle \frac{1}{\mathrm{k}}}{\displaystyle \frac{\theta }{t}}}`$ so that the observed frequencies are $`\omega =k\pm {\displaystyle \frac{\theta }{t}}{\displaystyle \frac{1}{2\mathrm{k}}}\left({\displaystyle \frac{\theta }{t^{}}}\right)^2+𝒪(\mathrm{k}^2)`$, that is, waves with opposite helicity propagate with different speeds. The results of astronomical observations lead to the constraint $`{\displaystyle \frac{\theta }{t}}10^{32}\text{eV}`$ .
Predicted effects of a hypothetical $`\theta `$ coupling on energy levels of atoms and existing atomic physics data imply the constraint $`\theta 10^8\text{m}^1`$ .
### 3.2 Test of well–posedness of Cauchy–problem
If for an evolution the Cauchy–problem is not well posed, then this evolution depends on its history or possesses memory. In mathematical terms this means, that, under certain circumstances, one has to pose all time derivatives up to infinite order: $`F`$, $`_tF`$, $`_t^2F`$, …, $`_t^nF`$, … . In a first approximation one may ask, whether there is any need to pose also $`_tF`$ in addition to $`F`$. For the usual Maxwell equations (here we assumne, for simplicity, a metric) we therefore have to add a term $`_t^2F`$:
$$\stackrel{~}{a}_\mu {}_{}{}^{\rho \sigma }_{0}^{2}F_{\rho \sigma }+_\nu F_\mu {}_{}{}^{\nu }=j_\mu .$$
(14)
The addition of such a term clearly violates LLI. In addition, if such a term is present, the homogeneous Maxwell equations reduce to constraints, because they are of first order only.
This term can be analyzed by the propagation of plane waves. The dispersion relation is
$`0`$ $`=`$ $`det\left(k^2\delta _{[\mu }^{[\rho }\delta _{\nu ]}^{\sigma ]}+\stackrel{~}{a}_{[\mu }{}_{}{}^{[\rho \sigma ]}k_{\nu ]}^{}\omega ^2\right)`$ (15)
$``$ $`(k^2)^5\left(\omega ^2\mathrm{k}^2+{\displaystyle \frac{1}{2}}\omega ^2\left(\stackrel{~}{a}_0{}_{}{}^{[0\widehat{\nu }]}k_{\widehat{\nu }]}^{}+\stackrel{~}{a}_{\widehat{\mu }}{}_{}{}^{[\widehat{\mu }0]}\omega +\stackrel{~}{a}_{\widehat{\mu }}{}_{}{}^{[\widehat{\mu }\widehat{\nu }]}k_{\widehat{\nu }}^{}\right)\right)`$
leading to the non–trivial solutions
$`\omega _{1,2}`$ $`=`$ $`\pm \mathrm{k}(1{\displaystyle \frac{i}{8}}(\stackrel{~}{a}_0{}_{}{}^{[0\widehat{\nu }]}+\stackrel{~}{a}_{\widehat{\mu }}{}_{}{}^{[\widehat{\mu }\widehat{\nu }]})k_{\widehat{\nu }}{\displaystyle \frac{i}{8}}\stackrel{~}{a}_{\widehat{\mu }}{}_{}{}^{[\widehat{\mu }0]}\mathrm{k})`$ (16)
$`\omega _3`$ $`=`$ $`{\displaystyle \frac{i}{\stackrel{~}{a}_{\widehat{\mu }}^{[\widehat{\mu }0]}}}(2+i{\displaystyle \frac{1}{2}}(\stackrel{~}{a}_0{}_{}{}^{[0\widehat{\nu }]}+\stackrel{~}{a}_{\widehat{\mu }}{}_{}{}^{[\widehat{\mu }\widehat{\nu }]})k_{\widehat{\nu }}+{\displaystyle \frac{1}{8}}(\stackrel{~}{a}_{\widehat{\mu }}{}_{}{}^{[\widehat{\mu }0]})^2\mathrm{k}^2).`$ (17)
The third solution has no physical relevance because it diverges for vanishing $`\stackrel{~}{a}`$. From the first two solutions we conclude that a second time–derivative in the Maxwell equations will lead to an anisotropic damping of plane waves.
### 3.3 Test of finite propagation speed
As for the heat or the Schrödinger equation, infinite propagation speed means that the order of spatial derivatives is larger than the the order of the time derivative. In our case this means that we have as modified Maxwell equations
$$_\nu F_\mu {}_{}{}^{\nu }+\stackrel{~}{b}_\mu {}_{}{}^{\widehat{\kappa }\widehat{\lambda }\rho \sigma }_{\widehat{\kappa }}^{}_{\widehat{\lambda }}F_{\rho \sigma }=j_\mu ,$$
(18)
which again violates LLI. Again, these equations can be analyzed through propagation phenomena. We get as dispersion relation
$`0`$ $`=`$ $`det\left(k^2\delta _{[\mu }^{[\rho }\delta _{\nu ]}^{\sigma ]}+k_{[\mu }\stackrel{~}{b}_{\nu ]}{}_{}{}^{\widehat{\kappa }\widehat{\lambda }[\rho \sigma ]}k_{\widehat{\kappa }}^{}k_{\widehat{\lambda }}\right)`$ (19)
$``$ $`(k^2)^5(\omega ^2\mathrm{k}^2+{\displaystyle \frac{1}{2}}(\stackrel{~}{b}_{\widehat{\mu }}{}_{}{}^{\widehat{\kappa }\widehat{\lambda }[\widehat{\mu }0]}\omega +(\stackrel{~}{b}_0{}_{}{}^{\widehat{\kappa }\widehat{\lambda }[0\widehat{\mu }]}+\stackrel{~}{b}_{\widehat{\nu }}{}_{}{}^{\widehat{\kappa }\widehat{\lambda }[\widehat{\mu }\widehat{\nu }]})k_{\widehat{\mu }})k_{\widehat{\kappa }}k_{\widehat{\lambda }})`$
with the solutions
$$\omega _{1,2}=\pm \mathrm{k}\frac{1}{4}\stackrel{~}{b}_{\widehat{\mu }}{}_{}{}^{\widehat{\kappa }\widehat{\lambda }[\widehat{\mu }0]}k_{\widehat{\kappa }}^{}k_{\widehat{\lambda }}\frac{1}{4\mathrm{k}}(\stackrel{~}{b}_0{}_{}{}^{\widehat{\kappa }\widehat{\lambda }[0\widehat{\mu }]}+\stackrel{~}{b}_{\widehat{\nu }}{}_{}{}^{\widehat{\kappa }\widehat{\lambda }[\widehat{\mu }\widehat{\nu }]})k_{\widehat{\mu }}k_{\widehat{\kappa }}k_{\widehat{\lambda }}$$
(20)
showing up an anisotropic dispersive propagation thus violating LLI.
## 4 Conclusion
We have presented a new test theory for the dynamics of the electromagnetic field. This theory was applied to propagation phenomena well suited to testing the Maxwell equations with astronomical data. We reviewed existing constraints on conventional anomalous electromagnetic field dynamics and also discussed the effects of anomalies associated with the appearance of higher-order field derivatives in the generalized Maxwell equations. These latter anomalies cause either anisotropic wave propagation or wave damping. The structure of the test theory reviewed here is sufficiently general to provide a context for the design and interpretation of experimental tests of special relativity and of metric gravitation theories like general relativity.
C.L. thanks the Optikzentrum of the University of Konstanz for financial support. |
warning/0002/hep-ph0002089.html | ar5iv | text | # Flavour-Conserving CP Phases in Supersymmetry and Implications for Exclusive 𝐵 Decays
## I Introduction
Within the standard model (SM), $`\mathrm{CP}`$ violation is caused by a non-zero complex phase in the Cabibbo-Kobayashi-Maskawa (CKM) quark mixing matrix . While the experimentally observed indirect $`\mathrm{CP}`$ violation in the neutral kaon system, $`ϵ_K`$, can be accommodated in the SM, it is still an open question whether the SM description of $`\mathrm{CP}`$ violation is consistent with the new experimental result on direct $`\mathrm{CP}`$ violation, $`ϵ^{}/ϵ_K`$, since the theoretical prediction of its precise value suffers from large hadronic uncertainties . On the other hand, if the baryon asymmetry of the universe has been generated via baryogenesis at the electroweak phase transition, the CKM mechanism of $`\mathrm{CP}`$ violation cannot account for the observed amount of baryon asymmetry. This feature could be a hint of the existence of $`\mathrm{CP}`$-violating sources outside the CKM matrix .
Important tests of the SM are provided by flavour-changing neutral current (FCNC) reactions involving $`B`$ decays , thus offering an opportunity to search for supersymmetric extensions of the SM . There are at present only a few FCNC processes which have been observed experimentally, but the situation will change considerably after the completion of $`B`$ factories in the near future.
In this work, we analyse the exclusive decays $`\overline{B}K^{()}l^+l^{}`$ and $`\overline{B}\rho (\pi )l^+l^{}`$ in the context of supersymmetry (SUSY) with minimal particle content and R-parity conservation . The inclusive reaction $`\overline{B}X_sl^+l^{}`$ within supersymmetric models has been extensively studied in Refs. and, more recently, in Ref. . New physics effects in the exclusive channels have been investigated in Refs. .
We place particular emphasis on $`\mathrm{CP}`$-violating effects associated with the partial rate asymmetry between $`B`$ and $`\overline{B}`$ decays as well as the forward-backward asymmetry of the $`l^{}`$. Within the SM these effects turn out to be unobservably small ($`10^3`$) in the decays $`\overline{B}K^{()}l^+l^{}`$ , and amount to only a few per cent in $`\overline{B}\rho (\pi )l^+l^{}`$ . However, in models with new $`\mathrm{CP}`$-violating phases in addition to the single phase of the CKM matrix, larger effects may occur due to the interference of amplitudes with different phases. The purpose of the present analysis is to explore $`\mathrm{CP}`$-violating observables in the aforementioned FCNC reactions that could provide evidence of a non-standard source of $`\mathrm{CP}`$ violation, and hence may be useful in analysing supersymmetry in future collider experiments.
The paper is organized as follows. In Sec. II, we exhibit the various mixing matrices of the minimal supersymmetric standard model (MSSM) in the presence of additional $`\mathrm{CP}`$-violating phases. Within such a framework we discuss different scenarios for the SUSY parameters. In Sec. III, we are primarily concerned with the short-distance matrix element and Wilson coefficients governing $`bs(d)l^+l^{}`$ in the MSSM. We also briefly describe an approximate procedure to incorporate quark antiquark resonant intermediate states – namely $`\rho ,\omega `$, and the $`J/\psi `$ family – which enter through the decay chain $`bs(d)V_{q\overline{q}}s(d)l^+l^{}`$. Section IV is devoted to the exclusive decay modes $`\overline{B}K^{()}l^+l^{}`$ and $`\overline{B}\pi (\rho )l^+l^{}`$, where formulae are given to calculate $`\mathrm{CP}`$ asymmetries which can be determined experimentally by measuring the difference of $`B`$ and $`\overline{B}`$ events. In Sec. V, we present our numerical results for $`\mathrm{CP}`$-violating observables in the non-resonant domain $`1.2\text{GeV}<M_{l^+l^{}}<M_{J/\psi }`$, taking into account experimental bounds on rare $`B`$ decays such as $`bs\gamma `$. We summarize and conclude in Sec. VI. The analytic formulae describing the short-distance effects in the presence of SUSY as well as the explicit expressions for the form factors are relegated to the Appendices.
## II The minimal supersymmetric standard model
In the MSSM, there are new sources of $`\mathrm{CP}`$ violation. In general, a large number of $`\mathrm{CP}`$-violating phases appear in the mass matrices as well as the couplings. After an appropriate redefinition of fields one ends up with at least two new $`\mathrm{CP}`$-violating phases, besides the phase of the CKM matrix and the QCD vacuum angle, which cannot be rotated away. For instance, in the MSSM with universal boundary conditions at some high scale only two new physical phases arise; namely $`\phi _{\mu _0}`$ associated with the Higgsino mass parameter $`\mu `$ in the superpotential and $`\phi _{A_0}`$ connected with the soft SUSY-breaking trilinear mass terms.
In order to fulfil the severe constraints on the electric dipole moments (EDM’s) of electron and neutron, one generally assumes that the new phases are less than $`𝒪(10^2)`$. Since there is no underlying symmetry which would force the phases to be small, this requires fine-tuning. Of course, one can relax the tight constraint on these phases by having masses of the superpartners in the TeV region; this heavy SUSY spectrum may, however, lead to an unacceptably large contribution to the cosmological relic density.
It has recently been pointed out by several authors that it is possible to evade the EDM constraints so that phases of $`𝒪(1)`$ still remain consistent with the current experimental upper limits. Methods that have been advocated to suppress the EDM’s include cancellations among different SUSY contributions , and nearly degenerate heavy sfermions for the first two generations while being consistent with naturalness bounds. The latter can be realized within the context of so-called ‘effective SUSY’ models , thereby solving the SUSY FCNC and $`\mathrm{CP}`$ problems.
To get an idea how supersymmetry affects $`\mathrm{CP}`$ observables in rare $`B`$ decays, we will consider as illustrative examples the following types of SUSY models:
* MSSM coupled to $`N=1`$ supergravity with a universal SUSY-breaking sector at the grand unification scale.
* Effective SUSY with near degeneracy of the heavy first and second generation sfermions.
In the present analysis, we restrict ourselves to the discussion of flavour-diagonal sfermion mass matrices – that is, we assume the CKM matrix to be the only source of flavour mixing.<sup>*</sup><sup>*</sup>*One should keep in mind that renormalization group effects induce flavour off-diagonal entries in the sfermion mass matrices at the weak scale (see below).
### A Mixing matrices and new $`\mathrm{𝐂𝐏}`$-violating phases
This subsection concerns the mass and mixing matrices relevant to our analysis. In what follows, we will adopt the conventions of Ref. .
#### 1 Charged Higgs-boson mass matrix
The mass-squared matrix of the charged Higgs bosons reads
$$M_{H^\pm }^2=\left(\begin{array}{cc}B\mu \mathrm{tan}\beta +M_W^2\mathrm{sin}^2\beta +t_1/v_1& B\mu +M_W^2\mathrm{sin}\beta \mathrm{cos}\beta \\ B\mu +M_W^2\mathrm{sin}\beta \mathrm{cos}\beta & B\mu \mathrm{cot}\beta +M_W^2\mathrm{cos}^2\beta +t_2/v_2\end{array}\right),$$
(1)
with
$$B\mu =\frac{1}{2}\mathrm{sin}2\beta (m_{H_1}^2+m_{H_2}^2+2|\mu |^2),$$
(2)
$$|\mu |^2=\frac{1}{2}M_Z^2+\frac{m_{H_2}^2\mathrm{sin}^2\beta m_{H_1}^2\mathrm{cos}^2\beta }{\mathrm{cos}2\beta }.$$
(3)
Here $`B`$ and $`\mu `$ refer to the complex soft SUSY-breaking and Higgsino mass parameters respectively, $`m_{H_{1,2}}^2`$ are the soft SUSY breaking Higgs-boson masses at the electroweak scale, and $`t_{1,2}`$ stand for the renormalized tadpoles. The mixing angle $`\beta `$ is defined as usual by $`\mathrm{tan}\beta v_2/v_1`$, with $`v_{1,2}`$ denoting the tree level vacuum expectation values (VEV’s) of the two neutral Higgs fields. In Eq. (2) we have adjusted the phase of the $`\mu `$ parameter in such a way that $`B\mu `$ is real at tree level, thereby ensuring that the VEV’s of the two Higgs fields are real. Consequently, the mass matrix becomes real and can be reduced to a diagonal form through a biorthogonal transformation $`(M_{H^\pm }^{\mathrm{diag}})^2=OM_{H^\pm }^2O^T`$. At the tree level, i.e. $`t_i=0`$ in Eq. (1), we have
$$O=\left(\begin{array}{cc}\mathrm{cos}\beta & \mathrm{sin}\beta \\ \mathrm{sin}\beta & \mathrm{cos}\beta \end{array}\right).$$
(4)
Before proceeding, we should mention that radiative corrections to the Higgs potential induce complex VEV’s. As a matter of fact, $`\mathrm{CP}`$ violation in the Higgs sector leads to an additional phase which, in the presence of chargino and neutralino contributions, cannot be rotated away by reparametrization of fields . As a result, the radiatively induced phase modifies the squark, chargino, and neutralino mass matrices. In the present analysis, we set this phase equal to zero.
#### 2 Squark mass matrices
We now turn to the $`6\times 6`$ squark mass-squared matrix which can be written as
$$M_{\stackrel{~}{q}}^2=\left(\begin{array}{cc}M_{\stackrel{~}{q}_{LL}}^2& M_{\stackrel{~}{q}_{LR}}^2e^{i\phi _{\stackrel{~}{q}}}\\ M_{\stackrel{~}{q}_{LR}}^2e^{i\phi _{\stackrel{~}{q}}}& M_{\stackrel{~}{q}_{RR}}^2\end{array}\right),\stackrel{~}{q}=\stackrel{~}{U},\stackrel{~}{D},$$
(5)
in the $`(\stackrel{~}{q}_L,\stackrel{~}{q}_R)`$ basis, and can be diagonalized by a unitary matrix $`R_{\stackrel{~}{q}}`$ such that
$$(M_{\stackrel{~}{q}}^{\mathrm{diag}})^2=R_{\stackrel{~}{q}}M_{\stackrel{~}{q}}^2R_{\stackrel{~}{q}}^{}.$$
(6)
For subsequent discussion it is useful to define the $`6\times 3`$ matrices
$$(\mathrm{\Gamma }^{q_L})_{ai}=(R_{\stackrel{~}{q}})_{ai},(\mathrm{\Gamma }^{q_R})_{ai}=(R_{\stackrel{~}{q}})_{a,i+3},q=U,D,$$
(7)
with $`U`$ and $`D`$ denoting up- and down-type quarks respectively. Working in the so-called ‘super-CKM’ basis in which the $`3\times 3`$ quark mass matrices $`M_U`$ and $`M_D`$ are real and diagonal, the submatrices in Eq. (5) take the form
$$M_{\stackrel{~}{U}_{LL}}^2=(M_{\stackrel{~}{U}}^2)_{LL}+M_U^2+\frac{1}{6}M_Z^2\mathrm{cos}2\beta (34\mathrm{sin}^2\theta _W)𝟙,$$
(9)
$$M_{\stackrel{~}{U}_{LR}}^2=M_U|A_U\mu ^{}\mathrm{cot}\beta 𝟙|,$$
(10)
$$M_{\stackrel{~}{U}_{RR}}^2=(M_{\stackrel{~}{U}}^2)_{RR}+M_U^2+\frac{2}{3}M_Z^2\mathrm{cos}2\beta \mathrm{sin}^2\theta _W𝟙,$$
(11)
$$\phi _{\stackrel{~}{U}}=\mathrm{arg}(A_U\mu ^{}\mathrm{cot}\beta 𝟙),$$
(12)
$$M_{\stackrel{~}{D}_{LL}}^2=V_{\mathrm{CKM}}^{}(M_{\stackrel{~}{U}}^2)_{LL}V_{\mathrm{CKM}}+M_D^2\frac{1}{6}M_Z^2\mathrm{cos}2\beta (32\mathrm{sin}^2\theta _W)𝟙,$$
(14)
$$M_{\stackrel{~}{D}_{LR}}^2=M_D|A_D\mu ^{}\mathrm{tan}\beta 𝟙|,$$
(15)
$$M_{\stackrel{~}{D}_{RR}}^2=(M_{\stackrel{~}{D}}^2)_{RR}+M_D^2\frac{1}{3}M_Z^2\mathrm{cos}2\beta \mathrm{sin}^2\theta _W𝟙,$$
(16)
$$\phi _{\stackrel{~}{D}}=\mathrm{arg}(A_D\mu ^{}\mathrm{tan}\beta 𝟙).$$
(17)
Here $`\theta _W`$ denotes the Weinberg angle, $`𝟙`$ represents a $`3\times 3`$ unit matrix, $`(M_{\stackrel{~}{q}}^2)_{LL}`$ and $`(M_{\stackrel{~}{q}}^2)_{RR}`$ are Hermitian scalar soft mass matrices, and $`V_{\mathrm{CKM}}`$ is the usual CKM matrix. In deriving Eq. (14), we have used the relation $`(M_{\stackrel{~}{D}}^2)_{LL}=V_{\mathrm{CKM}}^{}(M_{\stackrel{~}{U}}^2)_{LL}V_{\mathrm{CKM}}`$, which is due to SU(2) gauge invariance. Since we ignore flavour-mixing effects among squarks, $`(M_{\stackrel{~}{q}}^2)_{LL}`$ and $`(M_{\stackrel{~}{q}}^2)_{RR}`$ in Eqs. (7) and (II A 2) are diagonal – and hence real – whereas the $`A_q`$’s are given by
$$A_U=\mathrm{diag}(A_u,A_c,A_t),A_D=\mathrm{diag}(A_d,A_s,A_b),A_i|A_i|e^{i\phi _{A_i}}.$$
(18)
Consequently, the squark mass-squared matrix, Eq. (5), in the up-squark sector decomposes into a series of $`2\times 2`$ matrices. As far as the scalar top quark is concerned, we have
$$M_{\stackrel{~}{t}}^2=\left(\begin{array}{cc}m_{\stackrel{~}{t}_L}^2+m_t^2+\frac{1}{6}M_Z^2\mathrm{cos}2\beta (34\mathrm{sin}^2\theta _W)& m_t|A_t\mu ^{}\mathrm{cot}\beta |e^{i\phi _{\stackrel{~}{t}}}\\ m_t|A_t\mu ^{}\mathrm{cot}\beta |e^{i\phi _{\stackrel{~}{t}}}& m_{\stackrel{~}{t}_R}^2+m_t^2+\frac{2}{3}M_Z^2\mathrm{cos}2\beta \mathrm{sin}^2\theta _W\end{array}\right),$$
(19)
where $`m_{\stackrel{~}{t}_L}^2`$ and $`m_{\stackrel{~}{t}_R}^2`$ are diagonal elements of $`(M_{\stackrel{~}{U}}^2)_{LL}`$ and $`(M_{\stackrel{~}{U}}^2)_{RR}`$ respectively, while $`\phi _{\stackrel{~}{t}}`$ can be readily inferred from Eq. (12). Diagonalization of the stop mass-squared matrix then leads to the physical mass eigenstates $`\stackrel{~}{t}_1`$ and $`\stackrel{~}{t}_2`$, namely
$$\left(\begin{array}{c}\stackrel{~}{t}_1\\ \stackrel{~}{t}_2\end{array}\right)=\left(\begin{array}{cc}\mathrm{cos}\theta _{\stackrel{~}{t}}& \mathrm{sin}\theta _{\stackrel{~}{t}}e^{i\phi _{\stackrel{~}{t}}}\\ \mathrm{sin}\theta _{\stackrel{~}{t}}e^{i\phi _{\stackrel{~}{t}}}& \mathrm{cos}\theta _{\stackrel{~}{t}}\end{array}\right)\left(\begin{array}{c}\stackrel{~}{t}_L\\ \stackrel{~}{t}_R\end{array}\right)\left(\begin{array}{cc}\mathrm{\Gamma }_{33}^{U_L}& \mathrm{\Gamma }_{33}^{U_R}\\ \mathrm{\Gamma }_{63}^{U_L}& \mathrm{\Gamma }_{63}^{U_R}\end{array}\right)\left(\begin{array}{c}\stackrel{~}{t}_L\\ \stackrel{~}{t}_R\end{array}\right),$$
(20)
where the mixing angle $`\theta _{\stackrel{~}{t}}`$ is given by the expression ($`\pi /2\theta _{\stackrel{~}{t}}\pi /2`$)
$$\mathrm{tan}2\theta _{\stackrel{~}{t}}=\frac{2m_t|A_t\mu ^{}\mathrm{cot}\beta |}{(m_{\stackrel{~}{t}_L}^2m_{\stackrel{~}{t}_R}^2)+\frac{1}{6}M_Z^2\mathrm{cos}2\beta (38\mathrm{sin}^2\theta _W)}.$$
(21)
#### 3 Chargino mass matrix
The chargino mass matrix can be written as
$$M_{\stackrel{~}{\chi }^\pm }=\left(\begin{array}{cc}M_2& \sqrt{2}M_W\mathrm{sin}\beta \\ \sqrt{2}M_W\mathrm{cos}\beta & |\mu |e^{i\phi _\mu }\end{array}\right),$$
(22)
where we have adopted a phase convention in which the mass term of the W-ino field, $`M_2`$, is real and positive. Note that without loss of generality, we can always perform a global rotation to remove one of the three phases from the gaugino mass parameters $`M_i`$ ($`i=1,2,3`$).
The mass matrix can be cast in diagonal form by means of a biunitary transformation, namely
$$M_{\stackrel{~}{\chi }^\pm }^{\mathrm{diag}}=U^{}M_{\stackrel{~}{\chi }^\pm }V^{},$$
(23)
where $`M_{\stackrel{~}{\chi }^\pm }^{\mathrm{diag}}`$ is diagonal with positive eigenvalues, and $`U`$, $`V`$ are unitary matrices. Solving the eigenvalue problem
$$(M_{\stackrel{~}{\chi }^\pm }^{\mathrm{diag}})^2=U^{}M_{\stackrel{~}{\chi }^\pm }M_{\stackrel{~}{\chi }^\pm }^{}U^T=VM_{\stackrel{~}{\chi }^\pm }^{}M_{\stackrel{~}{\chi }^\pm }V^{},$$
(24)
we find
$$U=\left(\begin{array}{cc}\mathrm{cos}\theta _U& \mathrm{sin}\theta _Ue^{i\phi _U}\\ \mathrm{sin}\theta _Ue^{i\phi _U}& \mathrm{cos}\theta _U\end{array}\right),$$
(25)
$$V=\left(\begin{array}{cc}\mathrm{cos}\theta _Ve^{i\varphi _1}& \mathrm{sin}\theta _Ve^{i(\phi _V+\varphi _1)}\\ \mathrm{sin}\theta _Ve^{i(\phi _V\varphi _2)}& \mathrm{cos}\theta _Ve^{i\varphi _2}\end{array}\right),$$
(26)
with the mixing angles
$$\mathrm{tan}2\theta _U=\frac{2\sqrt{2}M_W[M_2^2\mathrm{cos}^2\beta +|\mu |^2\mathrm{sin}^2\beta +|\mu |M_2\mathrm{sin}2\beta \mathrm{cos}\phi _\mu ]^{1/2}}{M_2^2|\mu |^22M_W^2\mathrm{cos}2\beta },$$
(27)
$$\mathrm{tan}2\theta _V=\frac{2\sqrt{2}M_W[M_2^2\mathrm{sin}^2\beta +|\mu |^2\mathrm{cos}^2\beta +|\mu |M_2\mathrm{sin}2\beta \mathrm{cos}\phi _\mu ]^{1/2}}{M_2^2|\mu |^2+2M_W^2\mathrm{cos}2\beta },$$
(28)
$$\mathrm{tan}\phi _U=\frac{|\mu |\mathrm{sin}\phi _\mu \mathrm{sin}\beta }{M_2\mathrm{cos}\beta +|\mu |\mathrm{sin}\beta \mathrm{cos}\phi _\mu },$$
(29)
$$\mathrm{tan}\phi _V=\frac{|\mu |\mathrm{sin}\phi _\mu \mathrm{cos}\beta }{M_2\mathrm{sin}\beta +|\mu |\mathrm{cos}\beta \mathrm{cos}\phi _\mu },$$
(30)
$$\mathrm{tan}\varphi _1=\frac{|\mu |\mathrm{sin}\phi _\mu M_W^2\mathrm{sin}2\beta }{M_2(m_{\stackrel{~}{\chi }_1^\pm }^2|\mu |^2)+|\mu |M_W^2\mathrm{sin}2\beta \mathrm{cos}\phi _\mu },$$
(31)
$$\mathrm{tan}\varphi _2=\frac{|\mu |\mathrm{sin}\phi _\mu (m_{\stackrel{~}{\chi }_2^\pm }^2M_2^2)}{M_2M_W^2\mathrm{sin}2\beta +|\mu |(m_{\stackrel{~}{\chi }_2^\pm }^2M_2^2)\mathrm{cos}\phi _\mu }.$$
(32)
Here we have chosen $`\pi /2\theta _i\pi /2`$, $`\pi \phi _i,\varphi _i\pi `$, where $`i=U,V`$, and the chargino mass eigenvalues read
$`m_{\stackrel{~}{\chi }_{1,2}^\pm }^2={\displaystyle \frac{1}{2}}[M_2^2+|\mu |^2+2M_W^2`$ (33)
$``$ $`\{(M_2^2|\mu |^2)^2+4M_W^4\mathrm{cos}^22\beta +4M_W^2[M_2^2+|\mu |^2+2|\mu |M_2\mathrm{sin}2\beta \mathrm{cos}\phi _\mu ]\}^{1/2}].`$ (34)
### B SUSY particles and FCNC interactions
We present here the SUSY Lagrangian relevant to the FCNC processes of interest which will also serve as a means of fixing our notation. The interactions of charged Higgs bosons, charginos, neutralinos, and gluinos in the presence of new $`\mathrm{CP}`$ phases can be written as
$`_{\text{SUSY}}`$ $`=`$ $`{\displaystyle \frac{g}{\sqrt{2}M_W}}[\mathrm{cot}\beta (\overline{u}M_UV_{\mathrm{CKM}}P_Ld)+\mathrm{tan}\beta (\overline{u}V_{\mathrm{CKM}}M_DP_Rd)]H^+`$ (35)
$`+`$ $`{\displaystyle \underset{j=1}{\overset{2}{}}}\overline{\stackrel{~}{\chi }_j^{}}[\stackrel{~}{u}^{}(X_j^{U_L}P_L+X_j^{U_R}P_R)d+\stackrel{~}{\nu }^{}(X_j^{L_L}P_L+X_j^{L_R}P_R)l]`$ (36)
$`+`$ $`{\displaystyle \underset{k=1}{\overset{4}{}}}\overline{\stackrel{~}{\chi }_k^0}[\stackrel{~}{d}^{}(Z_k^{D_L}P_L+Z_k^{D_R}P_R)d+\stackrel{~}{l}^{}(Z_k^{L_L}P_L+Z_k^{L_R}P_R)l]`$ (37)
$``$ $`\sqrt{2}g_s{\displaystyle \underset{b=1}{\overset{8}{}}}\overline{\stackrel{~}{g}^b}\stackrel{~}{d}^{}(G^{D_L}P_LG^{D_R}P_R)T^bd+\mathrm{H}.\mathrm{c}.,`$ (38)
where generation indices have been suppressed, and $`P_{L,R}=(1\gamma _5)/2`$. The mixing matrices in the super-CKM basis are given by
$$X_j^{U_L}=g\left[V_{j1}^{}\mathrm{\Gamma }^{U_L}+V_{j2}^{}\mathrm{\Gamma }^{U_R}\frac{M_U}{\sqrt{2}M_W\mathrm{sin}\beta }\right]V_{\mathrm{CKM}},$$
(40)
$$X_j^{U_R}=gU_{j2}\mathrm{\Gamma }^{U_L}V_{\mathrm{CKM}}\frac{M_D}{\sqrt{2}M_W\mathrm{cos}\beta },$$
(41)
$$X_j^{L_L}=gV_{j1}^{}R_{\stackrel{~}{\nu }},X_j^{L_R}=gU_{j2}R_{\stackrel{~}{\nu }}\frac{M_E}{\sqrt{2}M_W\mathrm{cos}\beta },$$
(42)
$$Z_k^{D_L}=\frac{g}{\sqrt{2}}\left[\left(N_{k2}^{}+\frac{1}{3}\mathrm{tan}\theta _WN_{k1}^{}\right)\mathrm{\Gamma }^{D_L}+N_{k3}^{}\mathrm{\Gamma }^{D_R}\frac{M_D}{M_W\mathrm{cos}\beta }\right],$$
(43)
$$Z_k^{D_R}=\frac{g}{\sqrt{2}}\left[\frac{2}{3}\mathrm{tan}\theta _WN_{k1}\mathrm{\Gamma }^{D_R}+N_{k3}\mathrm{\Gamma }^{D_L}\frac{M_D}{M_W\mathrm{cos}\beta }\right],$$
(44)
$$Z_k^{L_L}=\frac{g}{\sqrt{2}}\left[(N_{k2}^{}+\mathrm{tan}\theta _WN_{k1}^{})\mathrm{\Gamma }^{L_L}N_{k3}^{}\mathrm{\Gamma }^{L_R}\frac{M_E}{M_W\mathrm{cos}\beta }\right],$$
(45)
$$Z_k^{L_R}=\frac{g}{\sqrt{2}}\left[2\mathrm{tan}\theta _WN_{k1}\mathrm{\Gamma }^{L_R}+N_{k3}\mathrm{\Gamma }^{L_L}\frac{M_E}{M_W\mathrm{cos}\beta }\right],$$
(46)
$$G^{D_L}=e^{i\phi _3/2}\mathrm{\Gamma }^{D_L},G^{D_R}=e^{i\phi _3/2}\mathrm{\Gamma }^{D_R},$$
(47)
$`\phi _3`$ being the phase of the gluino mass term $`M_3`$. (For scalar lepton as well as neutralino mass and mixing matrices, we refer the reader to Ref. .) In the remainder of this section, we briefly discuss two SUSY models with quite distinct scenarios for the $`\mathrm{CP}`$-violating phases.
### C Different scenarios for the SUSY parameters
#### 1 Constrained MSSM
In order to solve the FCNC problem in the MSSM and to further reduce the number of unknown parameters, the MSSM is generally embedded in a grand unified theory (GUT). This leads to the minimal supergravity (mSUGRA) inspired model, commonly referred to as the constrained MSSM . In this model one assumes universality of the soft terms at some high scale, which we take to be the scale of gauge coupling unification, $`M_{\mathrm{GUT}}`$, implying that (i) all gaugino mass parameters are equal to a common mass $`M_{1/2}`$; (ii) all the scalar mass parameters share a common value $`m_0`$; and (iii) all the soft trilinear couplings are equal to $`A_0`$. As a result, the mSUGRA model has only two new independent phases which are associated with the $`\mu _0`$ and $`A_0`$ parameters. After all, we have at the GUT scale
$$\mathrm{tan}\beta ,M_{1/2},m_0,|A_0|,|\mu _0|,\phi _{A_0},\phi _{\mu _0},$$
(48)
with $`M_{1/2}`$ and $`m_0`$ being real. The parameters at the electroweak scale are then obtained by solving the renormalization group equations (RGE’s).
A few remarks are in order here. First, the phases of the gaugino mass terms $`M_i`$ are not affected by the renormalization group evolution, and therefore the low energy gaugino mass parameters are real. Second, the phase that appears together with the $`\mu `$ parameter does not run at one-loop level so that $`\phi _\mu =\phi _{\mu _0}`$. Moreover, to satisfy the constraints on the EDM’s of electron and neutron, $`\phi _\mu `$ has to be of $`𝒪(10^2)`$ unless strong cancellations between different contributions occur. Third, solving the RGE’s for the evolution of the $`\mathrm{CP}`$-violating phase of the $`A_t`$-term yields a small value for $`\phi _{A_t}`$, and thus is not constrained by the EDM’s . Lastly, off-diagonal entries occur in the squark mass matrices due to renormalization group evolution of the parameters even in the absence of flavour mixing at the GUT scale. However, these effects are found to be small and therefore the squark mass matrix is essentially flavour diagonal at the electroweak scale (see also Refs. ).
#### 2 Effective supersymmetry
As an example of SUSY models with large CP phases, we consider the effective supersymmetry picture without assuming universality of sfermion masses at a high scale. Within such a framework, the first and second generation sfermions are almost degenerate and have masses above the TeV scale, while third generation sfermions can be light enough to be accessible at future hadron colliders. Consequently, FCNC reactions as well as one-loop contributions to the EDM’s of electron and neutron are well below the current experimental bounds.
However, it should be noted that the EDM’s also receive contributions at two-loop level involving scalar bottom and top quarks that may become important for phases of order unity in the large $`\mathrm{tan}\beta `$ regime . In our numerical work, $`\mathrm{tan}\beta `$ is assumed to be in the interval $`2\mathrm{tan}\beta 30`$.
## III Rare $`𝑩`$ decays and new physics
### A Short-distance matrix element
Let us start with the QCD-corrected matrix element describing the short-distance interactions in $`bs(d)l^+l^{}`$ within the SM. It is given by
$`_{\text{SD}}={\displaystyle \frac{G_F\alpha }{\sqrt{2}\pi }}V_{tb}V_{tf}^{}\{[(c_9^{\text{eff}}c_{10})H(k)|\overline{f}\gamma _\mu P_Lb|\overline{B}(p)`$ (50)
$`{\displaystyle \frac{2c_7^{\text{eff}}}{q^2}}H(k)|\overline{f}i\sigma _{\mu \nu }q^\nu (m_bP_R+m_fP_L)b|\overline{B}(p)]\overline{l}\gamma ^\mu P_Ll+(c_{10}c_{10})\overline{l}\gamma ^\mu P_Rl\},`$
where $`q`$ is the four-momentum of the lepton pair, and $`H=K,K^{}`$ ($`\pi ,\rho `$) in the case of $`f=s`$ ($`f=d`$). In the SM, the Wilson coefficients $`c_7^{\text{eff}}`$ and $`c_{10}`$ are real with values of $`0.314`$ and $`4.582`$ respectively, and the leading term in $`c_9^{\text{eff}}`$ has the form
$$c_9^{\text{eff}}=c_9+(3c_1+c_2)\left\{g(m_c,q^2)+\lambda _u\left[g(m_c,q^2)g(m_u,q^2)\right]\right\}+\mathrm{},$$
(51)
where $`c_9=4.216`$. The Wilson coefficients will be discussed in detail in Appendix B. In the above expression
$$\lambda _u\frac{V_{ub}V_{uf}^{}}{V_{tb}V_{tf}^{}}\{\begin{array}{c}\lambda ^2(\rho i\eta )\text{for}f=s,\hfill \\ \frac{\rho (1\rho )\eta ^2}{(1\rho )^2+\eta ^2}\frac{i\eta }{(1\rho )^2+\eta ^2}\text{for}f=d,\hfill \end{array}$$
(52)
with $`\rho `$ and $`\eta `$ being the Wolfenstein parameters , where the latter reflects the presence of $`\mathrm{CP}`$ violation in the SM. For definiteness, we will assume $`\rho =0.19`$ and $`\eta =0.35`$ .
Finally, the one-loop function $`g(m_i,q^2)`$ at the scale $`\mu _R=m_b`$ is given byIn order to avoid confusion with the $`\mu `$ parameter of the superpotential, we use the notation $`\mu _R`$ for the renormalization scale.
$`g(m_i,q^2)={\displaystyle \frac{8}{9}}\mathrm{ln}(m_i/m_b)+{\displaystyle \frac{8}{27}}+{\displaystyle \frac{4}{9}}y_i{\displaystyle \frac{2}{9}}(2+y_i)\sqrt{|1y_i|}`$ (54)
$`\times \left\{\mathrm{\Theta }(1y_i)\left[\mathrm{ln}\left({\displaystyle \frac{1+\sqrt{1y_i}}{1\sqrt{1y_i}}}\right)i\pi \right]+\mathrm{\Theta }(y_i1)2\mathrm{arctan}{\displaystyle \frac{1}{\sqrt{y_i1}}}\right\},`$
where $`y_i=4m_i^2/q^2`$. Observe that $`c\overline{c}`$ and $`u\overline{u}`$ loops provide absorptive parts that are mandatory, as we show below, for a non-zero partial width asymmetry besides the presence of a $`\mathrm{CP}`$-violating phase.
### B Wilson coefficients and new physics
Throughout this paper, we will assume that in the presence of new physics there are no new operators beyond those that correspond to the Wilson coefficients appearing in Eq. (50). (For a discussion of the implications of new operators for rare $`B`$ decays, see, e.g., Ref. .) Thus, the effect of new physics is simply to modify the matching conditions of the Wilson coefficients, i.e. their absolute values and phases at the electroweak scale.
As a result, we are left with additional SUSY contributions at one-loop level to the Wilson coefficients $`c_7^{\text{eff}}`$, $`c_9^{\text{eff}}`$, and $`c_{10}`$ in Eq. (50). In fact, they arise from penguin and box diagrams with (i) charged Higgs boson up-type quark loops; (ii) chargino up-type squark loops; (iii) neutralino down-type squark loops; and (iv) gluino down-type squark loops. Thus, the short-distance coefficients can be conveniently written as
$$c_i(M_W)=c_i^{\mathrm{SM}}(M_W)+c_i^{H^\pm }(M_W)+c_i^{\stackrel{~}{\chi }^\pm }(M_W)+c_i^{\stackrel{~}{\chi }^0}(M_W)+c_i^{\stackrel{~}{g}}(M_W)(i=7,\mathrm{},10).$$
(55)
The explicit expressions for the various Wilson coefficients are given in Appendix B. Since we limit our attention to flavour-conserving effects in the squark sector, the neutralino and gluino exchange contributions in Eq. (55) will be omitted in our numerical calculations.
For future reference, we parametrize the new physics contributions as follows:
$$R_i=\frac{c_i(M_W)}{c_i^{\mathrm{SM}}(M_W)}|R_i|e^{i\varphi _i},\chi =\frac{R_81}{R_71},$$
(56)
where $`\chi `$ is real to a good approximation within the models under study.
### C Resonant intermediate states
We have considered so far only the short-distance interactions. A more complete analysis, however, has also to take into account resonance contributions due to $`u\overline{u}`$, $`d\overline{d}`$, and $`c\overline{c}`$ intermediate states, i.e. $`\rho ,\omega ,J/\psi ,\psi ^{}`$, and so forth. A detailed discussion of the various theoretical suggestions of how to describe these effects is given in Ref. .
We employ here the approach proposed in Ref. which makes use of the renormalized photon vacuum polarization, $`\mathrm{\Pi }_{\mathrm{had}}^\gamma (s)`$, related to cross-section dataThis method assumes quark-hadron duality and rests on the factorization assumption.
$$R_{\mathrm{had}}(s)\frac{\sigma _{\text{tot}}(e^+e^{}\text{hadrons})}{\sigma (e^+e^{}\mu ^+\mu ^{})},$$
(57)
with $`sq^2`$. In fact, the absorptive part of the vacuum polarization is given by
$$\text{Im}\mathrm{\Pi }_{\mathrm{had}}^\gamma (s)=\frac{\alpha }{3}R_{\mathrm{had}}(s),$$
(58)
whereas the dispersive part may be obtained via a once-subtracted dispersion relation
$$\text{Re}\mathrm{\Pi }_{\mathrm{had}}^\gamma (s)=\frac{\alpha s}{3\pi }P\underset{4M_\pi ^2}{\overset{\mathrm{}}{}}\frac{R_{\mathrm{had}}(s^{})}{s^{}(s^{}s)}𝑑s^{},$$
(59)
with $`P`$ denoting the principal value. For example, in the case of the $`J/\psi `$ family (i.e. $`J/\psi ,\psi ^{},\mathrm{}`$) the imaginary part of the one-loop function $`g(m_c,s)`$, Eq. (54), can be expressed as
$$\text{Im}g(m_c,s)=\frac{\pi }{3}R_{\mathrm{had}}^{J/\psi }(s),R_{\mathrm{had}}^{J/\psi }(s)R_{\mathrm{cont}}^{c\overline{c}}(s)+R_{\mathrm{res}}^{J/\psi }(s),$$
(60)
where the subscripts ‘cont’ and ‘res’ stand for continuum and resonance contributions respectively, while the real part is given by
$$\text{Re}g(m_c,s)=\frac{8}{9}\mathrm{ln}(m_c/m_b)\frac{4}{9}+\frac{s}{3}P\underset{4M_\pi ^2}{\overset{\mathrm{}}{}}\frac{R_{\mathrm{had}}^{J/\psi }(s^{})}{s^{}(s^{}s)}𝑑s^{}.$$
(61)
The contributions from the continuum can be determined by means of experimental data given in Ref. , whereas the narrow resonances are well described by a relativistic Breit-Wigner distribution.
However, in order to reproduce correctly the branching ratio for direct $`J/\psi `$ production via the relation $`(V_{c\overline{c}}=J/\psi ,\psi ^{},\mathrm{}`$)
$$(BHV_{c\overline{c}}Hl^+l^{})=(BHV_{c\overline{c}})(V_{c\overline{c}}l^+l^{}),$$
(62)
where $`H`$ stands for pseudoscalar and scalar mesons, one has to multiply $`R_{\mathrm{res}}^{J/\psi }`$ in Eqs. (60) and (61) by a phenomenological factor $`\kappa `$, regardless of which method one uses for the description of the resonances .<sup>§</sup><sup>§</sup>§Strictly speaking, it is the term $`(3c_1+c_2)R_{\mathrm{res}}^{J/\psi }`$ – in the approximation of Eq. (51) – which has to be corrected to $`\kappa (3c_1+c_2)R_{\mathrm{res}}^{J/\psi }`$, taking into account non-factorizable contributions in two-body $`B`$ decays (see, e.g., Ref. ). Using the form factors of Ref. (see next section) together with experimental data on $`(BK^{()}J/\psi )`$, $`(BK^{()}\psi ^{})`$, and $`(B^{}\pi ^{}J/\psi )`$ , we find a magnitude for $`\kappa `$ of $`1.7`$ to $`3.3`$. At this point two remarks are in order. First, the branching ratio for direct $`J/\psi `$ and $`\psi ^{}`$ production is enhanced by a factor $`\kappa ^2`$, while it is essentially independent of $`\kappa `$ outside the resonance region. Second, the numerical results for average $`\mathrm{CP}`$ asymmetries in the non-resonant continuum $`1.2\text{GeV}<\sqrt{s}<2.9\text{GeV}`$ are not affected by the uncertainty in $`\kappa `$.
Similar considerations also hold for $`u\overline{u}`$ and $`d\overline{d}`$ systems except that the $`\rho `$ resonance is described through
$$R_{\mathrm{res}}^\rho (s)=\frac{1}{4}\left(1\frac{4M_\pi ^2}{s}\right)^{3/2}|F_\pi (s)|^2,$$
(63)
where the pion form factor is given by a modified Gounaris-Sakurai formula .
## IV The decays $`\overline{𝑩}\mathbf{}𝑲^{\mathbf{(}\mathbf{}\mathbf{)}}𝒍^\mathbf{+}𝒍^{\mathbf{}}`$ and $`\overline{𝑩}\mathbf{}𝝅\mathbf{(}𝝆\mathbf{)}𝒍^\mathbf{+}𝒍^{\mathbf{}}`$
The hadronic matrix elements in exclusive $`B`$ decays can be written in terms of $`q^2`$-dependent form factors, where $`q^2`$ is the invariant mass of the lepton pair. In the work described here, we employ heavy-to-light $`BK^{()}`$ and $`B\pi (\rho )`$ form factors determined by Melikhov and Nikitin within a relativistic quark model. To get an estimate of the theoretical uncertainty that is inherent to any model for the form factors, we also utilize the parametrization of Colangelo et al. , which makes use of QCD sum rule predictions.
For simplicity of presentation, we do not display corrections due to a non-zero lepton mass, which can be found in Refs. . (The same remark applies to the light quark masses $`m_{s,d}`$.) Henceforth we shall denote pseudoscalar and vector mesons by $`P=K,\pi `$ and $`V=K^{},\rho `$ respectively.
### A $`𝑩\mathbf{}𝑷`$ transitions
#### 1 Form factors
The hadronic matrix elements for the decays $`BP`$ can be parametrized in terms of three Lorentz-invariant form factors (see Appendix D for details), namely
$$P(k)|\overline{f}\gamma _\mu P_Lb|\overline{B}(p)=\frac{1}{2}[(2pq)_\mu f_+(q^2)+q_\mu f_{}(q^2)],$$
(65)
$$P(k)|\overline{f}i\sigma _{\mu \nu }q^\nu P_{L,R}b|\overline{B}(p)=\frac{1}{2}[(2pq)_\mu q^2(M_B^2M_P^2)q_\mu ]s(q^2),$$
(66)
with $`q=pk`$. Here we assume that the form factors are real, in the absence of final-state interactions. Note that the terms proportional to $`q_\mu `$ may be dropped in the case of massless leptons.
#### 2 Differential decay spectrum and forward-backward asymmetry
Introducing the shorthand notation
$$\lambda (a,b,c)=a^2+b^2+c^22(ab+bc+ac),$$
(67)
$$X_i=\frac{1}{2}\lambda ^{1/2}(M_B^2,M_i^2,s),$$
(68)
and recalling $`sq^2`$, the differential decay rate can be written as (neglecting $`m_l`$ and $`m_f`$)
$$\frac{d\mathrm{\Gamma }(\overline{B}Pl^+l^{})}{dsd\mathrm{cos}\theta _l}=\frac{G_F^2\alpha ^2}{2^8\pi ^5M_B^3}|V_{tb}V_{tf}^{}|^2X_P^3\left[|c_9^{\text{eff}}f_+(s)+2c_7^{\text{eff}}m_bs(s)|^2+|c_{10}f_+(s)|^2\right]\mathrm{sin}^2\theta _l.$$
(69)
Here $`\theta _l`$ is the angle between $`l^{}`$ and the outgoing hadron in the dilepton centre-of-mass system, and the Wilson coefficients are collected in Appendix B. Defining the forward-backward (FB) asymmetry as
$$A_{\text{FB}}(s)=\frac{{\displaystyle _0^1}d\mathrm{cos}\theta _l{\displaystyle \frac{d\mathrm{\Gamma }}{dsd\mathrm{cos}\theta _l}}{\displaystyle _1^0}d\mathrm{cos}\theta _l{\displaystyle \frac{d\mathrm{\Gamma }}{dsd\mathrm{cos}\theta _l}}}{{\displaystyle _0^1}d\mathrm{cos}\theta _l{\displaystyle \frac{d\mathrm{\Gamma }}{dsd\mathrm{cos}\theta _l}}+{\displaystyle _1^0}d\mathrm{cos}\theta _l{\displaystyle \frac{d\mathrm{\Gamma }}{dsd\mathrm{cos}\theta _l}}},$$
(70)
which is equivalent to the energy asymmetry discussed in Ref. , it follows directly from the distribution in Eq. (69) that $`A_{\text{FB}}`$ vanishes in the case of $`\overline{B}Pl^+l^{}`$ transitions. We note in passing that, given an extended operator basis (e.g. in models with neutral Higgs-boson exchange), new Dirac structures $`\overline{l}l`$ and $`\overline{l}\gamma _5l`$ may occur in Eq. (50), giving rise to a non-zero FB asymmetry in $`\overline{B}Pl^+l^{}`$.
### B $`𝑩\mathbf{}𝑽`$ transitions
#### 1 Form factors
The hadronic matrix elements describing the decays $`BV`$ are characterized by seven independent form factors, which we present in Appendix D, defined through ($`ϵ_{0123}=+1`$)
$`V(k)|\overline{f}\gamma _\mu P_Lb|\overline{B}(p)`$ (73)
$`=iϵ_{\mu \nu \alpha \beta }ϵ^\nu p^\alpha q^\beta g(q^2){\displaystyle \frac{1}{2}}\left\{ϵ_\mu ^{}f(q^2)+(ϵ^{}q)[(2pq)_\mu a_+(q^2)+q_\mu a_{}(q^2)]\right\},`$
$`V(k)|\overline{f}i\sigma _{\mu \nu }q^\nu P_{R,L}b|\overline{B}(p)=iϵ_{\mu \nu \alpha \beta }ϵ^\nu p^\alpha q^\beta g_+(q^2){\displaystyle \frac{1}{2}}ϵ_\mu ^{}[g_+(q^2)(M_B^2M_V^2)+q^2g_{}(q^2)]`$ (74)
$`\pm `$ $`{\displaystyle \frac{1}{2}}(ϵ^{}q)\left\{(2pq)_\mu \left[g_+(q^2)+{\displaystyle \frac{1}{2}}q^2h(q^2)\right]+q_\mu \left[g_{}(q^2){\displaystyle \frac{1}{2}}(M_B^2M_V^2)h(q^2)\right]\right\},`$ (75)
$`ϵ^\mu `$ being the polarization vector of the final-state meson, and $`q=pk`$.
#### 2 Differential decay spectrum and forward-backward asymmetry
The differential decay rate for $`\overline{B}Vl^+l^{}`$ in the case of massless leptons and light quarks takes the form ($`f=s`$ or $`d`$)
$$\frac{d\mathrm{\Gamma }(\overline{B}Vl^+l^{})}{dsd\mathrm{cos}\theta _l}=\frac{G_F^2\alpha ^2}{2^9\pi ^5M_B^3}|V_{tb}V_{tf}^{}|^2X_V\left[A(s)+B(s)\mathrm{cos}\theta _l+C(s)\mathrm{cos}^2\theta _l\right],$$
(76)
with $`X_V`$ as in Eq. (68). The quantities $`A`$, $`B`$, and $`C`$ are
$`A(s)`$ $`=`$ $`{\displaystyle \frac{2X_V^2}{M_V^2}}\left[sM_V^2\alpha _1(s)+{\displaystyle \frac{1}{4}}\left(1+{\displaystyle \frac{2sM_V^2}{X_V^2}}\right)\alpha _2(s)+X_V^2\alpha _3(s)+(kq)\alpha _4(s)\right],`$ (77)
$$B(s)=8X_V\text{Re}\{c_{10}^{}[c_9^{\text{eff}}sA_xA_yc_7^{\text{eff}}m_b(A_xB_y+A_yB_x)]\},$$
(78)
$$C(s)=\frac{2X_V^2}{M_V^2}\left[sM_V^2\alpha _1(s)\frac{1}{4}\alpha _2(s)X_V^2\alpha _3(s)(kq)\alpha _4(s)\right],$$
(79)
where $`kq=(M_B^2M_V^2s)/2`$, and
$$\alpha _1(s)=(|c_9^{\text{eff}}|^2+|c_{10}|^2)A_x^2+\frac{4|c_7^{\text{eff}}|^2m_b^2}{s^2}B_x^2\frac{4\text{Re}(c_7^{\text{eff}}c_9^{\text{eff}})m_b}{s}A_xB_x,$$
(81)
$$\alpha _2(s)=\alpha _1(s)_{xy},\alpha _3(s)=\alpha _1(s)_{xz},$$
(82)
$$\alpha _4(s)=(|c_9^{\text{eff}}|^2+|c_{10}|^2)A_yA_z+\frac{4|c_7^{\text{eff}}|^2m_b^2}{s^2}B_yB_z\frac{2\text{Re}(c_7^{\text{eff}}c_9^{\text{eff}})m_b}{s}(A_yB_z+A_zB_y),$$
(83)
in which the $`A_i`$’s and $`B_i`$’s are defined as
$$A_x=g(s),A_y=f(s),A_z=a_+(s),$$
(85)
$$B_x=g_+(s),B_y=g_+(s)(M_B^2M_V^2)+sg_{}(s),B_z=\left[g_+(s)+\frac{1}{2}sh(s)\right].$$
(86)
The complex Wilson coefficients $`c_7^{\text{eff}}`$, $`c_9^{\text{eff}}`$, and $`c_{10}`$ are given in Appendix B.
Finally, using Eqs. (70) and (76), we derive the forward-backward asymmetry
$$A_{\text{FB}}(s)=12X_V\frac{\text{Re}\{c_{10}^{}[c_9^{\text{eff}}sA_xA_yc_7^{\text{eff}}m_b(A_xB_y+A_yB_x)]\}}{[3A(s)+C(s)]}.$$
(87)
### C CP-violating observables
To discuss $`\mathrm{CP}`$-violating asymmetries, let us first recall the necessary ingredients. Suppose the decay amplitude for $`\overline{B}F`$ has the general form
$$𝒜(\overline{B}F)=e^{i\varphi _1}A_1e^{i\delta _1}+e^{i\varphi _2}A_2e^{i\delta _2},$$
(88)
where $`\delta _i`$ and $`\varphi _i`$ denote strong phases ($`\mathrm{CP}`$-conserving) and weak phases ($`\mathrm{CP}`$-violating) respectively ($`A_1`$ and $`A_2`$ being real). Together with the decay amplitude for the conjugate process
$$\overline{𝒜}(B\overline{F})=e^{i\varphi _1}A_1e^{i\delta _1}+e^{i\varphi _2}A_2e^{i\delta _2},$$
(89)
which can be obtained from Eq. (88) by means of $`\mathrm{CPT}`$ invariance, we may define the $`\mathrm{CP}`$ asymmetry as
$$A_{\mathrm{CP}}\frac{|𝒜|^2|\overline{𝒜}|^2}{|𝒜|^2+|\overline{𝒜}|^2}=\frac{2r\mathrm{sin}\varphi \mathrm{sin}\delta }{1+2r\mathrm{cos}\varphi \mathrm{cos}\delta +r^2},$$
(90)
with $`r=A_2/A_1`$, $`\varphi =\varphi _1\varphi _2`$, and $`\delta =\delta _1\delta _2`$. As can be easily seen from Eq. (90), a non-zero partial rate asymmetry requires the simultaneous presence of a $`\mathrm{CP}`$-violating phase $`\varphi `$ as well as a $`\mathrm{CP}`$-conserving dynamical phase $`\delta `$, the latter being provided by the one-loop functions $`g(m_c,s)`$ and $`g(m_u,s)`$ present in the Wilson coefficient $`c_9^{\text{eff}}`$ \[Eq. (51)\]. Notice that in the limit in which the charm quark mass equals the up quark mass there is no $`\mathrm{CP}`$ violation in the SM.
Given the differential decay distribution in the variables $`s`$ and $`\mathrm{cos}\theta _l`$, we can construct the following $`\mathrm{CP}`$-violating observables:
$$A_{\mathrm{CP}}^{D,S}(s)=\frac{{\displaystyle _{D,S}}d\mathrm{cos}\theta _l{\displaystyle \frac{d\mathrm{\Gamma }_{\text{diff}}}{dsd\mathrm{cos}\theta _l}}}{{\displaystyle _S}d\mathrm{cos}\theta _l{\displaystyle \frac{d\mathrm{\Gamma }_{\text{sum}}}{dsd\mathrm{cos}\theta _l}}},_{D,S}_0^1_1^0,$$
(92)
where we have introduced
$$\mathrm{\Gamma }_{\text{diff}}=\mathrm{\Gamma }(\overline{B}Hl^+l^{})\overline{\mathrm{\Gamma }}(B\overline{H}l^+l^{}),$$
(93)
$$\mathrm{\Gamma }_{\text{sum}}=\mathrm{\Gamma }(\overline{B}Hl^+l^{})+\overline{\mathrm{\Gamma }}(B\overline{H}l^+l^{}),$$
(94)
with $`H=K^{()},\pi ,\rho `$. It should be noted that the asymmetry $`A_{\mathrm{CP}}^D`$ represents a $`\mathrm{CP}`$-violating effect in the angular distribution of $`l^{}`$ in $`B`$ and $`\overline{B}`$ decays while $`A_{\mathrm{CP}}^S`$ is the asymmetry in the partial widths of these decays. As can be seen from Eqs. (69), (76), and (87), the latter involves the phases of $`c_7^{\text{eff}}`$ and $`c_9^{\text{eff}}`$ while the former is also sensitive to the phase of the Wilson coefficient $`c_{10}`$.
## V Numerical analysis
Given the SUSY contributions presented in the preceding sections, we now proceed to study the implications of supersymmetry for exclusive $`B`$ decays.
### A Experimental constraints
In our numerical analysis, we scan the SUSY parameter space as given in Ref. and take as input the parameters displayed in Appendix A. In addition, we take into account the following experimental constraints:
* From the measurement of the inclusive branching ratio $`(\overline{B}X_s\gamma )`$, which probes $`|c_7^{\text{eff}}|`$, one can derive upper and lower limits :
$$2.0\times 10^4<(\overline{B}X_s\gamma )<4.5\times 10^4(95\%\mathrm{C}.\mathrm{L}.).$$
(95)
This is specially useful to constrain extensions of the SM. Indeed, following the model-independent analysis performed in Ref. , and taking the Wilson coefficient $`c_7^{\text{eff}}`$ in leading-log approximation \[see Eq. (B2) of the Appendix\], we obtain
$`(\overline{B}X_s\gamma )`$ $``$ $`[0.801+0.444|R_7|^2+0.002|R_8|^2+1.192\text{Re}R_7`$ (97)
$`+0.083\text{Re}R_8+0.061\text{Re}(R_7R_8^{})]\times 10^4,`$
where $`R_7`$ and $`R_8`$ are as in Eq. (56). From Fig. 1 we infer that the present $`bs\gamma `$ measurement already excludes many solutions for $`R_7`$.
* A CDF search for the exclusive decays of interest yields the upper limits $`(B^0K^0\mu ^+\mu ^{})<4.0\times 10^6`$ and $`(B^+K^+\mu ^+\mu ^{})<5.2\times 10^6`$ at the $`90\%\mathrm{C}.\mathrm{L}.`$ . Note that the $`K^0\mu ^+\mu ^{}`$ upper limit is close to the SM prediction . As for the modes $`B\pi l^+l^{}`$ and $`B\rho l^+l^{}`$, we are not aware of any such limits.
* Non-observation of any SUSY signals at LEP 2 and the Tevatron imposes the following lower bounds :
$`m_{\stackrel{~}{\chi }^\pm }>86\text{GeV},m_{\stackrel{~}{\nu }}>43\text{GeV},m_{\stackrel{~}{t}_1}>86\text{GeV},m_{\stackrel{~}{q}}>260\text{GeV},m_{H^\pm }>90\text{GeV}.`$ (98)
### B $`\mathrm{𝐂𝐏}`$ asymmetries
As mentioned earlier, we investigate $`\mathrm{CP}`$ asymmetries in the low dilepton invariant mass region, i.e. $`1.2\text{GeV}<\sqrt{s}<(M_{J/\psi }200\text{MeV})`$, which is of particular interest because the low-$`s`$ region is sensitive to the Wilson coefficient $`c_7^{\text{eff}}`$ (in the case of $`BVl^+l^{}`$). In fact, it can receive large SUSY contributions and be complex, as well as change sign, while being consistent with the experimental measurement of $`bs\gamma `$. On the other hand, the new-physics effects are known to alter the remaining Wilson coefficients $`c_9^{\text{eff}}`$ and $`c_{10}`$ only slightly within mSUGRA and effective SUSY with no additional flavour structure beyond the usual CKM mechanism . Moreover, $`\mathrm{CP}`$ asymmetries above the $`J/\psi `$ resonance are dominated by $`c\overline{c}`$ resonant intermediate states, whereas below $`1\text{GeV}`$ the $`\rho `$ resonance has a strong influence on the asymmetry. This can be seen from Fig. 2,
where we show the $`\mathrm{CP}`$ asymmetry $`A_{\mathrm{CP}}^S`$ \[Eq. (92)\] between $`B^{}\pi ^{}l^+l^{}`$ and $`B^+\pi ^+l^+l^{}`$ as a function of the dilepton invariant mass within the SM and in the presence of SUSY contributions with new $`\mathrm{CP}`$-violating phases. It is evident that the predictions for $`\mathrm{CP}`$ asymmetries suffer from large theoretical uncertainties in the neighbourhood of the $`\rho `$ resonance and above the $`J/\psi `$, as discussed in Sec. III.
Using Eqs. (69) and (76) together with the definition for $`\mathrm{CP}`$-violating asymmetries, Eqs. (IV C), we can summarize our main findings as follows:
#### 1 CP violation in $`BPl^+l^{}`$
* The $`\mathrm{CP}`$-violating asymmetry in the $`l^{}`$ spectra of $`B`$ and $`\overline{B}`$ decays, $`A_{\mathrm{CP}}^D`$, vanishes in the case of $`BP`$ transitions. Within the framework of the constrained MSSM with phases of $`𝒪(10^2)`$ numerical values for the average asymmetry $`A_{\mathrm{CP}}^S`$ in the low-$`s`$ region are comparable to the SM predictions with asymmetries of $`0.1\%`$ and $`5\%`$ for $`bs`$ and $`bd`$ transitions respectively.
* Our results for $`A_{\mathrm{CP}}^S`$ between the decays $`\overline{B}Pl^+l^{}`$ and $`B\overline{P}l^+l^{}`$ in the context of effective SUSY with a light stop $`\stackrel{~}{t}_1`$ and phases of $`𝒪(1)`$ are shown in Figs. 3 and 4 for low and large $`\mathrm{tan}\beta `$ solutions which correspond to $`\text{Re}R_7>0`$ and $`\text{Re}R_7<0`$ respectively. Observe that the $`\mathrm{CP}`$-violating asymmetry $`A_{\mathrm{CP}}^S`$ in $`BP`$ depends only weakly on the sign and phase of $`c_7^{\text{eff}}`$. This is due to the fact that $`c_7^{\text{eff}}`$, which is constrained by the $`bs\gamma `$ measurement and not enhanced by a factor of $`1/s`$ in the low-$`s`$ region, is nearly one order of magnitude smaller than the leading term in $`c_9^{\text{eff}}`$ \[cf. Eq. (51)\].
* The $`\mathrm{CP}`$ asymmetry in the partial widths of $`\overline{B}Kl^+l^{}`$ and $`B\overline{K}l^+l^{}`$ changes sign for large values of $`\varphi _9`$, while $`|A_{\mathrm{CP}}^S|1\%`$ (see Fig. 3).
However, non-standard contributions to $`\varphi _9`$ are found to be small and hence $`A_{\mathrm{CP}}^S𝒪(10^3)`$. On the other hand, average asymmetries of $`(5`$$``$$`6)\%`$ are predicted for $`A_{\mathrm{CP}}^S`$ in the case of $`B\pi `$, even for values of $`\varphi _9`$ as small as $`10^2`$ (see Fig. 4).
Given a typical branching ratio of $`10^8`$ and a nominal asymmetry of $`6\%`$, a measurement at $`3\sigma `$ level requires $`3\times 10^{11}`$ $`b\overline{b}`$ pairs. (This rather challenging task might be feasible at LHC and the Tevatron.)
* The small magnitude of the $`\mathrm{CP}`$ asymmetry is also due to a suppression factor multiplying the indispensable absorptive part in $`c_9^{\text{eff}}`$, which is only slightly affected by new-physics contributions. Indeed, it follows from Eqs. (51) and (90) that
$$A_{\mathrm{CP}}r\mathrm{sin}\varphi \mathrm{sin}\delta \frac{(3c_1+c_2)}{|c_9|}\mathrm{sin}\varphi \mathrm{sin}\delta 10^2\mathrm{sin}\varphi \mathrm{sin}\delta ,$$
(99)
where the weak and strong phases $`\varphi `$ and $`\delta `$ can be of order unity.
* Numerical estimates for average $`\mathrm{CP}`$ asymmetries are mildly affected by the parametrization of form factors (see also Refs. ).
#### 2 CP violation in $`BVl^+l^{}`$
* The contribution of the Wilson coefficient $`c_7^{\text{eff}}`$ (or equivalently $`R_7`$) to the decay rate in $`BV`$ modes is enhanced by a factor of $`1/s`$ in the low-$`s`$ region. As seen from Fig. 5,
in the case of $`B\rho `$, the $`\mathrm{CP}`$-violating asymmetry $`A_{\mathrm{CP}}^S`$ can change sign for $`\mathrm{tan}\beta =2`$ (i.e. $`\text{Re}R_7>0`$), while for large $`\mathrm{tan}\beta `$ it is always negative. For small values of $`\varphi _9`$ an average $`\mathrm{CP}`$ asymmetry of about $`5\%`$ is predicted for both $`\mathrm{tan}\beta =2`$ and $`\mathrm{tan}\beta =30`$ solutions. Since the distributions of $`A_{\mathrm{CP}}^S`$ for $`BK^{}`$ are very similar to the ones obtained for $`BK`$, we refrain from showing the corresponding plots.
* As we have already mentioned, the $`\mathrm{CP}`$-violating asymmetry in the angular distribution of $`l^{}`$ in $`B`$ and $`\overline{B}`$ decays can, in principle, probe the phase of the Wilson coefficient $`c_{10}`$. This is shown in Fig. 6,
where we have plotted the $`\mathrm{CP}`$ asymmetry as a function of $`\varphi _9`$ and $`\varphi _{10}`$ for large $`\mathrm{tan}\beta `$ (i.e. $`\text{Re}R_7<0`$). Unfortunately, the mSUGRA and effective SUSY predictions for the average asymmetry $`A_{\mathrm{CP}}^D`$ turn out to be unobservably small.
## VI Discussion and conclusions
In this paper, we have studied the consequences of new $`\mathrm{CP}`$-violating phases for exclusive $`B`$ decays within the framework of supersymmetric extensions of the SM, ignoring intergenerational mixing in the squark sector. We have examined $`\mathrm{CP}`$-violating asymmetries in the partial widths as well as angular distribution of $`l^{}`$ between the exclusive channels $`bs(d)l^+l^{}`$ and $`\overline{b}\overline{s}(\overline{d})l^+l^{}`$ in the invariant mass region $`1.2\text{GeV}<M_{l^+l^{}}<2.9\text{GeV}`$. The essential conclusion of our analysis is that it is not sufficient to have additional $`\mathrm{CP}`$ phases of $`𝒪(1)`$ in order to obtain large $`\mathrm{CP}`$-violating effects.
Within the constrained MSSM and effective SUSY with a complex CKM matrix and additional $`\mathrm{CP}`$ phases, we obtain values for the average asymmetry $`A_{\mathrm{CP}}^S`$ of about $`6\%`$ ($`5\%`$) in the decays $`\overline{B}\pi (\rho )l^+l^{}`$ and $`B\overline{\pi }(\overline{\rho })l^+l^{}`$, taking into account experimental constraints on EDM’s of electron and neutron, as well as rare $`B`$ decays such as $`bs\gamma `$. As for the asymmetry in the angular distribution, $`A_{\mathrm{CP}}^D`$, it probes the phase of the Wilson coefficient $`c_{10}`$, but will be unobservable at future colliders. Numerical estimates of the $`\mathrm{CP}`$ asymmetries in the decays $`\overline{B}K^{()}l^+l^{}`$ and $`B\overline{K}^{()}l^+l^{}`$ turn out to be small (less than $`1\%`$) and are comparable to the SM result.
Our analysis shows that the smallness in $`\mathrm{CP}`$ asymmetries is mainly due to the coefficient $`(3c_1+c_2)/|c_9|`$ which multiplies the requisite absorptive part in $`c_9^{\text{eff}}`$ \[Eq. (51)\], and which is only slightly affected by the new-physics contributions discussed in Sec. III. Therefore, any sizable $`\mathrm{CP}`$-violating effect in the low-$`s`$ region requires large non-standard contributions to the short-distance coefficient $`c_9^{\text{eff}}`$ and/or $`c_{10}`$, as well as additional $`\mathrm{CP}`$ phases of $`𝒪(1)`$. By the same token, any large $`\mathrm{CP}`$-violating effect would provide a clue to physics beyond the SM. A detailed discussion of this point will be given elsewhere.
One could argue, however, that the inclusion of flavour off-diagonal contributions (i.e. gluino and neutralino diagrams) to the Wilson coefficients might lead to higher $`\mathrm{CP}`$ asymmetries. In fact, it has been pointed out in Refs. that even in the presence of large supersymmetric $`\mathrm{CP}`$ phases, a non-trivial flavour structure in the soft-breaking terms is necessary in order to obtain sizable contributions to $`\mathrm{CP}`$ violation in the $`K`$ system and to $`\mathrm{CP}`$ asymmetries in two-body neutral $`B`$ decays (see also Ref. ). Using the mass insertion approximation, such effects have recently been studied in Ref. which predicts a partial width asymmetry for $`bsl^+l^{}`$ of a few per cent in the low-$`s`$ domain.
Finally, we would like to recall that, nevertheless, large $`\mathrm{CP}`$ asymmetries may occur in rare $`B`$ decays like the observed $`bs\gamma `$ modes, where $`A_{\mathrm{CP}}`$ can be substantial (up to $`\pm 45\%`$) in some part of the parameter space .
###### Acknowledgements.
F.K. would like to thank Gudrun Hiller for useful discussions. This research has been supported by the TMR Network of the EC under contract ERBFMRX-CT96-0090.
## A Numerical inputs
Unless otherwise specified, we use the experimental values as compiled by the Particle Data Group and the parameters displayed in Eq. (A3).
$`\begin{array}{c}m_t=175\text{GeV},m_b=4.8\text{GeV},m_c=1.4\text{GeV},m_s=170\text{MeV},m_d=10\text{MeV},\hfill \\ m_u=5\text{MeV},\alpha =1/129,\mathrm{\Lambda }_{\text{QCD}}=225\text{MeV}.\hfill \end{array}`$ (A3)
## B Wilson coefficients and SUSY
For the sake of convenience, we provide in this appendix formulae for the Wilson coefficients $`c_7^{\text{eff}}`$, $`c_9^{\text{eff}}`$, and $`c_{10}`$ in the presence of SUSY, using the results derived in Refs. . Since we study the case of massless leptons, we retain only those contributions that do not vanish in the limit $`m_l0`$. As for $`\tau `$ leptons in the final state, there are further charged and neutral Higgs-boson contributions \[see also Eqs. (II B)\].
Introducing the shorthand notation
$$\eta _s=\frac{\alpha _s(M_W)}{\alpha _s(m_b)},r_W=\frac{m_t^2}{M_W^2},r_{H^\pm }=\frac{m_t^2}{m_{H^\pm }^2},r_B^A=\frac{m_A^2}{m_B^2},$$
(B1)
the Wilson coefficient $`c_7^{\text{eff}}`$ evaluated at $`\mu _R=m_b`$ has the form (in leading-log approximation)
$$c_7^{\text{eff}}=\eta _s^{16/23}c_7(M_W)+\frac{8}{3}(\eta _s^{14/23}\eta _s^{16/23})c_8(M_W)+\underset{i=1}{\overset{8}{}}h_i\eta _s^{a_i},$$
(B2)
with the coefficients $`a_i`$, $`h_i`$ tabulated in Ref. . Recalling Eqs. (II B) and (55), and using the one-loop functions $`f_i`$ listed in Appendix C, the various contributions to $`c_{7,8}(M_W)`$ can be written as follows:
* Standard model:
$$c_7^{\mathrm{SM}}(M_W)=\frac{1}{4}r_Wf_1(r_W).$$
(B3)
* Charged Higgs boson:
$$c_7^{H^\pm }(M_W)=\frac{1}{12}[r_{H^\pm }f_1(r_{H^\pm })\mathrm{cot}^2\beta +2f_2(r_{H^\pm })].$$
(B4)
* Chargino:Notice that the one-loop function appearing in the last term of Eq. (B5) is actually $`f_2+5/2`$. However, using the explicit form for the squark mixing matrices \[Eqs. (II B)\], the constant term vanishes – reflecting the unitarity of the mixing matrices.
$`c_7^{\stackrel{~}{\chi }^\pm }(M_W)={\displaystyle \frac{1}{6g^2V_{tb}V_{tf}^{}}}{\displaystyle \underset{a=1}{\overset{6}{}}}{\displaystyle \underset{j=1}{\overset{2}{}}}{\displaystyle \frac{M_W^2}{m_{\stackrel{~}{\chi }_j^\pm }^2}}`$ (B5)
$`\times `$ $`\left[(X_j^{U_L})_{na}(X_j^{U_L})_{a3}f_1(r_{\stackrel{~}{\chi }_j^\pm }^{\stackrel{~}{u}_a})2(X_j^{U_L})_{na}(X_j^{U_R})_{a3}{\displaystyle \frac{m_{\stackrel{~}{\chi }_j^\pm }}{m_b}}f_2(r_{\stackrel{~}{\chi }_j^\pm }^{\stackrel{~}{u}_a})\right],`$ (B6)
where we have defined
$$n=\{\begin{array}{c}1\text{for}f=d,\hfill \\ 2\text{for}f=s.\hfill \end{array}$$
(B7)
For completeness, we also give the expressions for the neutralino and gluino contributions which vanish in the limit of flavour-diagonal squark mass matrices.
* Neutralino:
$`c_7^{\stackrel{~}{\chi }^0}(M_W)`$ (B8)
$`=`$ $`{\displaystyle \frac{1}{6g^2V_{tb}V_{tf}^{}}}{\displaystyle \underset{a=1}{\overset{6}{}}}{\displaystyle \underset{k=1}{\overset{4}{}}}{\displaystyle \frac{M_W^2}{m_{\stackrel{~}{\chi }_k^0}^2}}\left[(Z_k^{D_L})_{na}(Z_k^{D_L})_{a3}f_3(r_{\stackrel{~}{\chi }_k^0}^{\stackrel{~}{d}_a})+2(Z_k^{D_L})_{na}(Z_k^{D_R})_{a3}{\displaystyle \frac{m_{\stackrel{~}{\chi }_k^0}}{m_b}}f_4(r_{\stackrel{~}{\chi }_k^0}^{\stackrel{~}{d}_a})\right].`$ (B9)
* Gluino:
$$c_7^{\stackrel{~}{g}}(M_W)=\frac{4g_s^2}{9g^2V_{tb}V_{tf}^{}}\underset{a=1}{\overset{6}{}}\frac{M_W^2}{m_{\stackrel{~}{g}}^2}\left[(G^{D_L})_{na}(G^{D_L})_{a3}f_3(r_{\stackrel{~}{g}}^{\stackrel{~}{d}_a})2(G^{D_L})_{na}(G^{D_R})_{a3}\frac{m_{\stackrel{~}{g}}}{m_b}f_4(r_{\stackrel{~}{g}}^{\stackrel{~}{d}_a})\right].$$
(B10)
The corresponding expressions $`c_8^{\mathrm{SM}}(M_W),\mathrm{},c_8^{\stackrel{~}{\chi }^0}(M_W)`$ are obtained changing $`f_ig_i`$ in Eqs. (B3)–(B8), with $`g_i`$ collected in Appendix C, while the gluino contribution reads
$$c_8^{\stackrel{~}{g}}(M_W)=\frac{4g_s^2}{9g^2V_{tb}V_{tf}^{}}\underset{a=1}{\overset{6}{}}\frac{M_W^2}{m_{\stackrel{~}{g}}^2}\left[(G^{D_L})_{na}(G^{D_L})_{a3}g_5(r_{\stackrel{~}{g}}^{\stackrel{~}{d}_a})2(G^{D_L})_{na}(G^{D_R})_{a3}\frac{m_{\stackrel{~}{g}}}{m_b}g_6(r_{\stackrel{~}{g}}^{\stackrel{~}{d}_a})\right].$$
(B11)
The Wilson coefficient $`c_9^{\text{eff}}`$ at $`\mu _R=m_b`$ in next-to-leading approximation is given by
$`c_9^{\text{eff}}`$ $`=`$ $`c_9\left[1+{\displaystyle \frac{\alpha _s(m_b)}{\pi }}\omega (s/m_b^2)\right]+g(m_c,s)\left(3c_1+c_2+3c_3+c_4+3c_5+c_6\right)`$ (B12)
$`+`$ $`\lambda _u\left[g(m_c,s)g(m_u,s)\right]\left(3c_1+c_2\right){\displaystyle \frac{1}{2}}g(m_f,s)\left(c_3+3c_4\right)`$ (B13)
$``$ $`{\displaystyle \frac{1}{2}}g(m_b,s)\left(4c_3+4c_4+3c_5+c_6\right)+{\displaystyle \frac{2}{9}}\left(3c_3+c_4+3c_5+c_6\right),`$ (B14)
where $`\lambda _u`$ and $`g(m_i,s)`$ are defined in Eqs. (52) and (54) respectively, with $`sq^2`$. As far as the Wilson coefficients $`c_1`$$`c_6`$ are concerned, we have numerically
$$c_1=0.249,c_2=1.108,c_3=0.011,c_4=0.026,c_5=0.007,c_6=0.031,$$
(B15)
using the values given in Appendix A. Further,
$$c_9=c_9(M_W)\frac{4}{9}+P_0+P_E\underset{i}{}E^i,$$
(B16)
with $`i=\mathrm{SM},H^\pm ,\stackrel{~}{\chi }^\pm ,\stackrel{~}{\chi }^0,\stackrel{~}{g}`$, and
$$c_9(M_W)=\underset{i}{}\left(\frac{Y^i}{\mathrm{sin}^2\theta _W}4Z^i\right)+\frac{4}{9},$$
(B17)
where the analytic expressions for $`P_0`$, $`P_E`$, and $`E^{\mathrm{SM}}`$ are given in Ref. . Since $`P_EP_0`$, we shall keep only the SM contribution in the last term of Eq. (B16). Moreover, as discussed in Ref. , the order $`\alpha _s`$ correction in Eq. (B12) due to one-gluon exchange may be regarded as a contribution to the form factors, and hence we set $`\omega =0`$ in Eq. (B12).
Turning to the Wilson coefficient $`c_{10}`$, it has the simple form
$$c_{10}(M_W)=\underset{i}{}\frac{Y^i}{\mathrm{sin}^2\theta _W}.$$
(B18)
Note that the corresponding operator does not renormalize and thus $`c_{10}(M_W)=c_{10}(\mu _R)`$.
The expressions for the various contributions to $`Y`$ and $`Z`$ read as follows:Regarding the expressions for the chargino and neutralino box-diagram contributions, and the sign discrepancy between Ref. and Ref. , we confirm the results of the latter.
* Standard model:
$$Y^{\mathrm{SM}}=\frac{1}{8}f_9(r_W),Z^{\mathrm{SM}}=\frac{1}{72}f_{10}(r_W).$$
(B19)
* Charged Higgs:
$$Y^{H^\pm }Y_Z^{H^\pm },Z^{H^\pm }Z_\gamma ^{H^\pm }+Z_Z^{H^\pm },$$
(B21)
$$Z_\gamma ^{H^\pm }=\frac{1}{72}f_6(r_{H^\pm })\mathrm{cot}^2\beta ,Y_Z^{H^\pm }=Z_Z^{H^\pm }=\frac{1}{8}r_Wf_5(r_{H^\pm })\mathrm{cot}^2\beta .$$
(B22)
* Chargino:
$$Y^{\stackrel{~}{\chi }^\pm }Y_Z^{\stackrel{~}{\chi }^\pm }+Y_{\text{box}}^{\stackrel{~}{\chi }^\pm },Z^{\stackrel{~}{\chi }^\pm }Z_\gamma ^{\stackrel{~}{\chi }^\pm }+Z_Z^{\stackrel{~}{\chi }^\pm },$$
(B24)
$$Z_\gamma ^{\stackrel{~}{\chi }^\pm }=\frac{1}{36g^2V_{tb}V_{tf}^{}}\underset{a=1}{\overset{6}{}}\underset{j=1}{\overset{2}{}}\frac{M_W^2}{m_{\stackrel{~}{u}_a}^2}[(X_j^{U_L})_{na}(X_j^{U_L})_{a3}f_7(r_{\stackrel{~}{u}_a}^{\stackrel{~}{\chi }_j^\pm })],$$
(B25)
$`Y_Z^{\stackrel{~}{\chi }^\pm }=Z_Z^{\stackrel{~}{\chi }^\pm }={\displaystyle \frac{1}{2g^2V_{tb}V_{tf}^{}}}{\displaystyle \underset{a,b=1}{\overset{6}{}}}{\displaystyle \underset{i,j=1}{\overset{2}{}}}\{(X_i^{U_L})_{na}(X_j^{U_L})_{b3}[c_2(m_{\stackrel{~}{\chi }_i^\pm }^2,m_{\stackrel{~}{u}_a}^2,m_{\stackrel{~}{u}_b}^2)(\mathrm{\Gamma }^{U_L}\mathrm{\Gamma }^{U_L})_{ab}\delta _{ij}`$ (B26)
$``$ $`c_2(m_{\stackrel{~}{u}_a}^2,m_{\stackrel{~}{\chi }_i^\pm }^2,m_{\stackrel{~}{\chi }_j^\pm }^2)\delta _{ab}V_{i1}^{}V_{j1}+{\displaystyle \frac{1}{2}}m_{\stackrel{~}{\chi }_i^\pm }m_{\stackrel{~}{\chi }_j^\pm }c_0(m_{\stackrel{~}{u}_a}^2,m_{\stackrel{~}{\chi }_i^\pm }^2,m_{\stackrel{~}{\chi }_j^\pm }^2)\delta _{ab}U_{i1}U_{j1}^{}]\},`$ (B27)
$$Y_{\text{box}}^{\stackrel{~}{\chi }^\pm }=\frac{M_W^2}{g^2V_{tb}V_{tf}^{}}\underset{a=1}{\overset{6}{}}\underset{i,j=1}{\overset{2}{}}[(X_i^{U_L})_{na}(X_j^{U_L})_{a3}d_2(m_{\stackrel{~}{\chi }_i^\pm }^2,m_{\stackrel{~}{\chi }_j^\pm }^2,m_{\stackrel{~}{u}_a}^2,m_{\stackrel{~}{\nu }_{1,2}}^2)V_{i1}^{}V_{j1}],$$
(B28)
with $`m_{\stackrel{~}{\nu }_1}(m_{\stackrel{~}{\nu }_2})`$ in the case of $`e^+e^{}(\mu ^+\mu ^{})`$ in the final state.
* Neutralino:
$$Y^{\stackrel{~}{\chi }^0}Y_Z^{\stackrel{~}{\chi }^0}+Y_{\text{box}}^{\stackrel{~}{\chi }^0},Z^{\stackrel{~}{\chi }^0}Z_\gamma ^{\stackrel{~}{\chi }^0}+Z_Z^{\stackrel{~}{\chi }^0}+Z_{\text{box}}^{\stackrel{~}{\chi }^0},$$
(B30)
$$Z_\gamma ^{\stackrel{~}{\chi }^0}=\frac{1}{216g^2V_{tb}V_{tf}^{}}\underset{a=1}{\overset{6}{}}\underset{k=1}{\overset{4}{}}\frac{M_W^2}{m_{\stackrel{~}{d}_a}^2}[(Z_k^{D_L})_{na}(Z_k^{D_L})_{a3}f_8(r_{\stackrel{~}{d}_a}^{\stackrel{~}{\chi }_k^0})],$$
(B31)
$`Y_Z^{\stackrel{~}{\chi }^0}=Z_Z^{\stackrel{~}{\chi }^0}={\displaystyle \frac{1}{2g^2V_{tb}V_{tf}^{}}}{\displaystyle \underset{a,b=1}{\overset{6}{}}}{\displaystyle \underset{k,l=1}{\overset{4}{}}}\{(Z_k^{D_L})_{na}(Z_l^{D_L})_{b3}`$ (B32)
$`\times `$ $`[c_2(m_{\stackrel{~}{\chi }_k^0}^2,m_{\stackrel{~}{d}_a}^2,m_{\stackrel{~}{d}_b}^2)(\mathrm{\Gamma }^{D_R}\mathrm{\Gamma }^{D_R})_{ab}\delta _{kl}c_2(m_{\stackrel{~}{d}_a}^2,m_{\stackrel{~}{\chi }_k^0}^2,m_{\stackrel{~}{\chi }_l^0}^2)\delta _{ab}(N_{k3}^{}N_{l3}N_{k4}^{}N_{l4})`$ (B34)
$`{\displaystyle \frac{1}{2}}m_{\stackrel{~}{\chi }_k^0}m_{\stackrel{~}{\chi }_l^0}c_0(m_{\stackrel{~}{d}_a}^2,m_{\stackrel{~}{\chi }_k^0}^2,m_{\stackrel{~}{\chi }_l^0}^2)\delta _{ab}(N_{k3}N_{l3}^{}N_{k4}N_{l4}^{})]\},`$
$`Y_{\text{box}}^{\stackrel{~}{\chi }^0}=2\mathrm{sin}^2\theta _WZ_{\text{box}}^{\stackrel{~}{\chi }^0}+{\displaystyle \frac{M_W^2}{2g^2V_{tb}V_{tf}^{}}}{\displaystyle \underset{a=1}{\overset{6}{}}}{\displaystyle \underset{k,l=1}{\overset{4}{}}}\{(Z_k^{D_L})_{na}(Z_l^{D_L})_{a3}`$ (B35)
$`\times `$ $`[d_2(m_{\stackrel{~}{\chi }_k^0}^2,m_{\stackrel{~}{\chi }_l^0}^2,m_{\stackrel{~}{d}_a}^2,m_{\stackrel{~}{l}_{1,2}}^2)(N_{k2}^{}+\mathrm{tan}\theta _WN_{k1}^{})(N_{l2}+\mathrm{tan}\theta _WN_{l1})`$ (B37)
$`+{\displaystyle \frac{1}{2}}m_{\stackrel{~}{\chi }_k^0}m_{\stackrel{~}{\chi }_l^0}d_0(m_{\stackrel{~}{\chi }_k^0}^2,m_{\stackrel{~}{\chi }_l^0}^2,m_{\stackrel{~}{d}_a}^2,m_{\stackrel{~}{l}_{1,2}}^2)(N_{k2}+\mathrm{tan}\theta _WN_{k1})(N_{l2}^{}+\mathrm{tan}\theta _WN_{l1}^{})]\},`$
$`Z_{\text{box}}^{\stackrel{~}{\chi }^0}={\displaystyle \frac{M_W^2}{g^2V_{tb}V_{tf}^{}}}{\displaystyle \underset{a=1}{\overset{6}{}}}{\displaystyle \underset{k,l=1}{\overset{4}{}}}\{(Z_k^{D_L})_{na}(Z_l^{D_L})_{a3}\mathrm{sec}^2\theta _W`$ (B38)
$`\times `$ $`[d_2(m_{\stackrel{~}{\chi }_k^0}^2,m_{\stackrel{~}{\chi }_l^0}^2,m_{\stackrel{~}{d}_a}^2,m_{\stackrel{~}{l}_{4,5}}^2)N_{k1}^{}N_{l1}+{\displaystyle \frac{1}{2}}m_{\stackrel{~}{\chi }_k^0}m_{\stackrel{~}{\chi }_l^0}d_0(m_{\stackrel{~}{\chi }_k^0}^2,m_{\stackrel{~}{\chi }_l^0}^2,m_{\stackrel{~}{d}_a}^2,m_{\stackrel{~}{l}_{4,5}}^2)N_{k1}N_{l1}^{}]\},`$ (B39)
with $`m_{\stackrel{~}{l}_{1,4}}(m_{\stackrel{~}{l}_{2,5}})`$ for $`e^+e^{}(\mu ^+\mu ^{})`$ in the final state.
* Gluino:
$$Y^{\stackrel{~}{g}}Y_Z^{\stackrel{~}{g}},Z^{\stackrel{~}{g}}Z_\gamma ^{\stackrel{~}{g}}+Z_Z^{\stackrel{~}{g}},$$
(B41)
$$Z_\gamma ^{\stackrel{~}{g}}=\frac{g_s^2}{81g^2V_{tb}V_{tf}^{}}\underset{a=1}{\overset{6}{}}\frac{M_W^2}{m_{\stackrel{~}{d}_a}^2}[(G^{D_L})_{na}(G^{D_L})_{a3}f_8(r_{\stackrel{~}{d}_a}^{\stackrel{~}{g}})],$$
(B42)
$$Y_Z^{\stackrel{~}{g}}=Z_Z^{\stackrel{~}{g}}=\frac{4g_s^2}{3g^2V_{tb}V_{tf}^{}}\underset{a,b=1}{\overset{6}{}}[(G^{D_L})_{na}(G^{D_L})_{b3}c_2(m_{\stackrel{~}{g}}^2,m_{\stackrel{~}{d}_a}^2,m_{\stackrel{~}{d}_b}^2)(G^{D_R}G^{D_R})_{ab}].$$
(B43)
The functions $`f_i`$, $`c_i`$, and $`d_i`$ are given in Eqs. (C)–(C) below.
## C Auxiliary functions
Here we list the functions $`f_i`$, $`g_i`$, $`c_i`$, and $`d_i`$ introduced in the previous section.
$$f_1(x)=\frac{7+5x+8x^2}{6(1x)^3}\frac{x(23x)}{(1x)^4}\mathrm{ln}x,$$
(C2)
$$f_2(x)=\frac{x(35x)}{2(1x)^2}+\frac{x(23x)}{(1x)^3}\mathrm{ln}x,$$
(C3)
$$f_3(x)=\frac{2+5xx^2}{6(1x)^3}+\frac{x}{(1x)^4}\mathrm{ln}x,$$
(C4)
$$f_4(x)=\frac{1+x}{2(1x)^2}+\frac{x}{(1x)^3}\mathrm{ln}x,$$
(C5)
$$f_5(x)=\frac{x}{1x}+\frac{x}{(1x)^2}\mathrm{ln}x,$$
(C6)
$$f_6(x)=\frac{x(3879x+47x^2)}{6(1x)^3}+\frac{x(46x+3x^3)}{(1x)^4}\mathrm{ln}x,$$
(C7)
$$f_7(x)=\frac{52101x+43x^2}{6(1x)^3}+\frac{69x+2x^3}{(1x)^4}\mathrm{ln}x,$$
(C8)
$$f_8(x)=\frac{27x+11x^2}{(1x)^3}+\frac{6x^3}{(1x)^4}\mathrm{ln}x,$$
(C9)
$$f_9(x)=\frac{x(4x)}{1x}+\frac{3x^2}{(1x)^2}\mathrm{ln}x,$$
(C10)
$$f_{10}(x)=\frac{x(108259x+163x^218x^3)}{2(1x)^3}\frac{850x+63x^2+6x^324x^4}{(1x)^4}\mathrm{ln}x,$$
(C11)
$$g_1(x)=3f_3(x),g_2(x)=3[f_4(x)1/2],g_3(x)=3f_3(x),g_4(x)=3f_4(x),$$
(C13)
$$g_5(x)=\frac{1}{8}\left[\frac{1140x19x^2}{2(1x)^3}+\frac{3x(19x)}{(1x)^4}\mathrm{ln}x\right],$$
(C14)
$$g_6(x)=\frac{3}{8}\left[\frac{513x}{(1x)^2}+\frac{x(19x)}{(1x)^3}\mathrm{ln}x\right],$$
(C15)
$$c_0(m_1^2,m_2^2,m_3^2)=[\frac{m_1^2\mathrm{ln}(m_1^2/\mu _R^2)}{(m_1^2m_2^2)(m_1^2m_3^2)}+(m_1m_2)+(m_1m_3)],$$
(C17)
$$c_2(m_1^2,m_2^2,m_3^2)=\frac{3}{8}\frac{1}{4}[\frac{m_1^4\mathrm{ln}(m_1^2/\mu _R^2)}{(m_1^2m_2^2)(m_1^2m_3^2)}+(m_1m_2)+(m_1m_3)],$$
(C18)
$`d_0(m_1^2,m_2^2,m_3^2,m_4^2)`$ (C20)
$`=`$ $`[{\displaystyle \frac{m_1^2\mathrm{ln}(m_1^2/\mu _R^2)}{(m_1^2m_2^2)(m_1^2m_3^2)(m_1^2m_4^2)}}+(m_1m_2)+(m_1m_3)+(m_1m_4)],`$ (C21)
$`d_2(m_1^2,m_2^2,m_3^2,m_4^2)`$ (C22)
$`=`$ $`{\displaystyle \frac{1}{4}}[{\displaystyle \frac{m_1^4\mathrm{ln}(m_1^2/\mu _R^2)}{(m_1^2m_2^2)(m_1^2m_3^2)(m_1^2m_4^2)}}+(m_1m_2)+(m_1m_3)+(m_1m_4)].`$ (C23)
## D Form factors
In Tables I and II we summarize the two different sets of form factors discussed in Sec. IV, which are related via
$$F_1(q^2)=f_+(q^2),$$
(D2)
$$F_0(q^2)=f_+(q^2)+\frac{q^2}{M_B^2M_P^2}f_{}(q^2),$$
(D3)
$$F_T(q^2)=(M_B+M_P)s(q^2),$$
(D4)
$$V(q^2)=(M_B+M_V)g(q^2),$$
(D5)
$$A_1(q^2)=\frac{f(q^2)}{M_B+M_V},$$
(D6)
$$A_2(q^2)=(M_B+M_V)a_+(q^2),$$
(D7)
$$A_0(q^2)=\frac{q^2a_{}(q^2)+f(q^2)+(M_B^2M_V^2)a_+(q^2)}{2M_V},$$
(D8)
$$T_1(q^2)=\frac{1}{2}g_+(q^2),$$
(D9)
$$T_2(q^2)=\frac{1}{2}\left[g_+(q^2)+\frac{q^2g_{}(q^2)}{M_B^2M_V^2}\right],$$
(D10)
$$T_3(q^2)=\frac{1}{2}\left[g_{}(q^2)\frac{(M_B^2M_V^2)h(q^2)}{2}\right].$$
(D11) |
warning/0002/cond-mat0002132.html | ar5iv | text | # Long-time-tail Effects on Lyapunov Exponents of a Random, Two-dimensional Field-driven Lorentz Gas
## 1 Introduction
The Lorentz gas has proved to be a useful model for studying the relations between dynamical systems theory and non-equilibrium properties of many body systems. This model consists of a set of scatterers that are fixed in space together with moving particles that collide with the scatterers. Here we consider the version of the model in two dimensions where the scatterers are fixed hard disks, placed at random in the plane without overlapping. Each of the moving particles is a point particle with a mass and a charge, and is subjected to an external, uniform electric field as well as a Gaussian thermostat which is designed to keep the kinetic energy of the moving particle at a constant value. The particles make elastic, specular collisions with the scatterers, but do not interact with each other. The interest in the Lorentz gas model stems from the fact that its chaotic properties can be analyzed in some detail, at least if the scatterers form a sufficiently dilute, quenched gas, so that the average distance between scatterers is large compared to their radii. The interest in a thermostatted electric field arises from the fact that at small fields a transport coefficient, the electrical conductivity of the particles, is proportional to the sum of the Lyapunov exponents describing the chaotic motion of the moving particle . The Lyapunov exponents are to be calculated for the case where the charged particle is described by a non-equilibrium steady state phase-space distribution function which is reached from some typical initial distribution function after a sufficiently long period of time. In this state, the distribution function for an ensemble of moving particles (all interacting with the scatterers and the field, but not with each other) is independent of time and its average over the distribution of scatterers is spatially homogeneous. It is known from computer simulations and theoretical discussions that in the stationary state the trajectories of the moving particles in phase space lie on a fractal attractor of lower dimension than the dimension of the constant energy surface, which is three dimensional for the constant energy Lorentz gas in two dimensions. There can be at most two non-zero Lyapunov exponents for this model since the Lyapunov exponent in the direction of the phase-space trajectory is zero. Also, the relation between the Lyapunov exponents and the electrical conductivity requires that the sum of the non-zero exponents should be negative due to the positivity of electrical conductivity .
The case of the dilute, random Lorentz gas has already been studied in detail. Van Beijeren and coworkers have calculated the Lyapunov spectrum for an equilibrium Lorentz gas in two and three dimensions using various kinetic theory methods including Boltzmann equation techniques. These methods were also applied to the dilute, random Lorentz gas in a thermostatted electric field with results for the Lyapunov exponents that are in excellent agreement with computer simulations . Moreover, the results for the field-dependent case were in accord with the relation between the electrical conductivity and the Lyapunov exponents for the moving particle.
The purpose of this paper is to extend the results obtained for the Lyapunov exponents for dilute Lorentz gases to higher densities. Our central themes will be: (a) to describe a general method, based upon the BBGKY hierarchy equations, for accomplishing this task, and (b) to examine the effects on the Lyapunov exponents of long range in time correlations between the moving particle and the scatterers produced by correlated collision sequences where the particle collides with a given scatterer more than once and the time interval between such re-collisions is on the order of several mean free times, with an arbitrary number of intermediate collisions with other scatterers. These correlated collision sequences are of particular interest in kinetic and transport theory because they are responsible for the “long-time-tail” effects in the Green-Kubo time correlation functions, which lead to various divergences in the transport coefficients for two and three dimensional gases, where all of the particles move . In the case of a Lorentz gas in $`d`$ dimensions, the Green-Kubo velocity correlation functions decay with time, $`t`$, as $`t^{(d/2+1)}`$ and the diffusion coefficient is finite in both two and three dimensions. Here we describe the effects of these type of correlations on the Lyapunov exponents for the two-dimensional Lorentz gas, in equilibrium, where we find no effect, and in a thermostatted electric field, where we find a small, logarithmic dependence of the Lyapunov exponents upon the applied field. This logarithmic effect is an indicator for similar effects to be expected when one calculates Lyapunov exponents associated with more general transport in two-dimensional gases, whereas in three dimensional systems one would expect corresponding non-analytic terms proportional to $`\stackrel{~}{ϵ}^{\mathrm{\hspace{0.17em}5}/2}`$. In the case of the Lorentz gas, at least, the logarithmic terms can easily be associated with the logarithmic terms that appear in the field dependent collision frequency, and a very simple argument can be used to establish this relation between logarithmic terms in the Lyapunov exponents and in the collision frequency.
In Section 2 of this article we describe the general theory of Lyapunov exponents of a two-dimensional thermostatted electric field-driven Lorentz gas and quote the results within the scope of the Boltzmann equation. In Section 3, we generalize the theory to incorporate the effect of correlated collision sequences. In Section 4, we outline the calculation of the effects of the correlated collision sequences on the non-zero Lyapunov exponents using the BBGKY equations discussed in Section 3 and obtain the non-analytic field-dependent term in the Lyapunov exponents, originating from the correlated collision sequences, along with other analytic field-dependent terms. In Section 5, we present some simple arguments explaining the field-dependence of the collision frequency and show that this is the sole origin of the non-analytic, field-dependent terms in the Lyapunov exponents. Notice that the arguments given in Section 5 are independent of and much simpler to follow than the formalism developed in Sections 3 and 4. We conclude in Section 6 with a discussion of the results obtained here, and with a consideration of open questions. Methods for determination of the field-dependence of the collision frequency are outlined in the Appendix.
## 2 Lyapunov exponents of field driven Lorentz gases in two dimensions
### 2.1 General theory
The random Lorentz gas consists of point particles of mass $`m`$ and charge $`q`$ moving in a random array of fixed scatterers. In two dimensions, each scatterer is a hard disk of radius $`a`$. The disks do not overlap with each other and are distributed with number density $`n`$, such that at low density $`na^2<1`$. The point particles are acted upon by a uniform, constant electric field $`\stackrel{}{E}`$ in the $`\widehat{x}`$ direction, but there is no interaction between any two point particles. There is also a Gaussian thermostat in the system to keep the speed of each particle constant at $`v`$ by means of a dynamical friction during flights between collisions with the scatterers. The collisions between a point particle and the scatterers are instantaneous, specular and elastic. During a flight, the equations of motion of a point particle are
$`\dot{\stackrel{}{r}}=\stackrel{}{v}={\displaystyle \frac{\stackrel{}{p}}{m}},\dot{\stackrel{}{p}}=m\dot{\stackrel{}{v}}=q\stackrel{}{E}\alpha \stackrel{}{p}`$ (1)
Fig. 1 : Collision between a point particle and a scatterer.
and at a collision with a scatterer, the post-collisional position and velocity, $`\stackrel{}{r}_+`$ and $`\stackrel{}{v}_+`$, are related to the pre-collisional position and velocity, $`\stackrel{}{r}_{}`$ and $`\stackrel{}{v}_{}`$, by
$`\stackrel{}{r}_+=\stackrel{}{r}_{},\stackrel{}{v}_+=\stackrel{}{v}_{}\mathrm{\hspace{0.17em}2}(\stackrel{}{v}_{}\widehat{\sigma })\widehat{\sigma },`$ (2)
where $`\widehat{\sigma }`$ is the unit vector from the center of the scatterer to the point of collision (see Fig. 1). The fact that each particle has a constant speed $`v`$ determines the value of $`\alpha `$ :
$`\alpha ={\displaystyle \frac{q\stackrel{}{E}\stackrel{}{p}}{p^2}}\dot{\stackrel{}{p}}=q\stackrel{}{E}{\displaystyle \frac{q\stackrel{}{E}\stackrel{}{p}}{p^2}}\stackrel{}{p}.`$ (3)
Equivalently, in polar coordinates, the velocity direction with respect to the field, defined through $`\widehat{v}\widehat{x}=\mathrm{cos}\theta `$, changes between collisions as
$`\dot{\theta }=\epsilon \mathrm{sin}\theta ,`$ (4)
where $`\epsilon ={\displaystyle \frac{q|\stackrel{}{E}|}{mv}}`$ and we define the dimensionless electric field parameter $`\stackrel{~}{\epsilon }={\displaystyle \frac{\epsilon l}{v}}`$, where $`l=(2na)^{\mathrm{\hspace{0.17em}1}}`$ is the mean free path length for the particle in the dilute Lorentz gas. To denote the electric field, we will normally use $`\epsilon `$, though from time to time we will use $`\stackrel{~}{\epsilon }`$, too.
Treating this two-dimensional Lorentz gas as a dynamical system, we define the Lyapunov exponents in the usual way: a point particle in its phase space $`(\stackrel{}{r},\stackrel{}{v})=\stackrel{}{𝐗}`$ starts at time $`t_0`$ at a phase space location $`\stackrel{}{𝐗}(t_0)`$. Under time evolution, $`\stackrel{}{𝐗}(t)`$ follows a trajectory in this phase space which we call the “reference trajectory”. We consider an infinitesimally displaced trajectory which starts at the same time $`t_0`$, but at $`\stackrel{}{𝐗^{\frac{}{}}}(t_0)=\stackrel{}{𝐗}(t_0)+\delta \stackrel{}{𝐗}(t_0)`$. Under time evolution, $`\stackrel{}{𝐗^{\frac{}{}}}(t)`$ follows another trajectory, always staying infinitesimally close to the reference trajectory. This trajectory we call the “adjacent trajectory”. Typically the two trajectories will separate in time due to the convex nature of the collisions. Thus, we can define the positive Lyapunov exponent as
$`\lambda _+=\underset{\frac{}{T\mathrm{}}}{lim}\underset{\frac{}{|\delta \stackrel{}{𝐗}(t_0)|0}}{lim}{\displaystyle \frac{1}{T}}\mathrm{ln}{\displaystyle \frac{|\delta \stackrel{}{𝐗}(t_0+T)|}{|\delta \stackrel{}{𝐗}(t_0)|}}.`$ (5)
for a typical trajectory of the system.
We assume that, for small fields, this Lorentz gas system is hyperbolic. Since the two-dimensional Lorentz gas can have at most two nonzero Lyapunov exponents, we denote the negative Lyapunov exponent by $`\lambda _{}`$. Without any loss of generality, we can choose to measure the separation of the reference and adjacent trajectories equivalently in $`\stackrel{}{r}`$-space, thereby reducing the definition of the positive Lyapunov exponent to
$`\lambda _+=\underset{\frac{}{T\mathrm{}}}{lim}\underset{\frac{}{|\delta \stackrel{}{r}(t_0)|0}}{lim}{\displaystyle \frac{1}{T}}\mathrm{ln}{\displaystyle \frac{|\delta \stackrel{}{r}(t_0+T)|}{|\delta \stackrel{}{r}(t_0)|}}.`$ (6)
In order to calculate the right hand side of Eq. (6), we introduce another dynamical quantity, the radius of curvature $`\rho `$, characterizing the spatial separation of the two trajectories (see Fig. 2) :
Fig. 2 : $`\rho (t)={\displaystyle \frac{\delta S(t)}{\delta \theta (t)}}=|AP|`$.
In Fig. 2, a particle on the reference trajectory would be at point A at time $`t`$. At the same time, a particle on the adjacent trajectory would be at B. A local perpendicular on the reference trajectory at A intersects the adjacent trajectory at C. The backward extensions of instantaneous velocity directions on the reference and adjacent trajectories at A and C, respectively, intersect each other at point P. We denote the length of the line segment AC by $`\delta S(t)`$ and $`\mathrm{}APC`$ by $`\delta \theta (t)`$. The radius of curvature associated with the particle on the reference trajectory at time $`t`$ is then given by
$`\rho (t)={\displaystyle \frac{\delta S(t)}{\delta \theta (t)}}=|AP|.`$ (7)
Having defined $`\rho (t)`$, one can make a simple geometric argument to show that
$$\delta \dot{S}(t)=v\delta \theta (t)=\frac{v\delta S(t)}{\rho (t)},$$
(8)
so as to obtain a version of Sinai’s formula ,
$`\lambda =\underset{T\mathrm{}}{lim}{\displaystyle \frac{v}{T}}{\displaystyle _{\frac{}{t_0}}^{\frac{}{t_0+T}}}{\displaystyle \frac{dt}{\rho (t)}}.`$ (9)
During a flight, the equation of motion for $`\rho `$ is given by
$`\dot{\rho }=v+\rho \epsilon \mathrm{cos}\theta +{\displaystyle \frac{\rho ^2\epsilon ^2\mathrm{sin}^2\theta }{v}}.`$ (10)
At a collision with a scatterer, the post-collisional velocity angle $`\theta _+`$ and radius of curvature $`\rho _+`$ are related to the pre-collisional velocity angle $`\theta _{}`$ and radius of curvature $`\rho _{}`$ by :
$`\theta _+=\theta _{}\pi +2\varphi ,{\displaystyle \frac{1}{\rho _+}}={\displaystyle \frac{1}{\rho _{}}}+{\displaystyle \frac{2}{a\mathrm{cos}\varphi }}+{\displaystyle \frac{\epsilon }{v}}\mathrm{tan}\varphi (\mathrm{sin}\theta _{}+\mathrm{sin}\theta _+),`$ (11)
where $`\varphi `$ is the collision angle, i.e, $`\mathrm{cos}\varphi =|\widehat{v}_{}\widehat{\sigma }|=|\widehat{v}_+\widehat{\sigma }|`$ (see Fig. 1).
Now we assume that, for sufficiently weak electric field, the field-driven Lorentz gas in two dimensions is ergodic, and that we can replace the long time average in Eq. (6) by a non-equilibrium steady state (NESS) average, including an average over all allowed configurations of scatterers, to obtain
$`\lambda =<{\displaystyle \frac{v}{\rho }}>_{\text{NESS}}.`$ (12)
The electric field is considered weak if the work done by the electric field on the point particle over a flight of one mean free path is much smaller than the particle’s kinetic energy, i.e, $`{\displaystyle \frac{q|\stackrel{}{E}|l}{mv^2}}={\displaystyle \frac{\epsilon l}{v}}=\stackrel{~}{\epsilon }<<\mathrm{\hspace{0.17em}1}`$.
We note for future reference, that the sum of the two nonzero Lyapunov exponents is related to the average of the friction coefficient $`\alpha `$, by
$$\lambda _++\lambda _{}=<\alpha >_{\text{NESS}}=<\frac{q\stackrel{}{E}\stackrel{}{v}}{mv^2}>_{\text{NESS}}=\frac{\stackrel{}{J}\stackrel{}{E}}{mv^2}=\frac{\sigma E^2}{mv^2}.$$
(13)
Here the electric current $`\stackrel{}{J}=q\stackrel{}{v}_{\text{NESS}}`$ is, for small fields, assumed to satisfy Ohm’s law, $`\stackrel{}{J}=\sigma \stackrel{}{E}`$, and $`\sigma `$ is the electrical conductivity.
### 2.2 Results obtained using the Lorentz-Boltzmann equation
To the lowest order in density, one can assume that the collisions suffered by the point particle are uncorrelated, and use an extended Lorentz-Boltzmann equation (ELBE) for the distribution function of the moving particle, $`f_1(\stackrel{}{r},\stackrel{}{v},\rho ,t)`$ in $`(\stackrel{}{r},\stackrel{}{v},\rho )`$-space needed for the evaluation of the averages appearing in Eqs. (12) and (13). To calculate the positive Lyapunov exponent, one needs to consider the forward-time ELBE while to calculate the negative Lyapunov exponent one needs the time reversed ELBE. To the leading order in density, the Lyapunov exponents are then given by :
$`\lambda _+^{\text{(B)}}=\lambda _0{\displaystyle \frac{11}{48}}{\displaystyle \frac{l}{v}}\epsilon ^2+O(\epsilon ^4)\text{and}\lambda _{}^{\text{(B)}}=\lambda _0{\displaystyle \frac{7}{48}}{\displaystyle \frac{l}{v}}\epsilon ^2+O(\epsilon ^4).`$ (14)
The superscript, B, indicates that these are results obtained from the Lorentz-Boltzmann equation. Here $`\lambda _0`$ is the positive Lyapunov exponent for a field-free Lorentz gas (see for example ) given by
$$\lambda _0=2nav[\mathrm{\hspace{0.17em}1}𝒞\mathrm{ln}(2na^2)],$$
(15)
where $`𝒞`$ is Euler’s constant, $`𝒞=0.5772\mathrm{}`$. From Eqs. (13) and (14), using Einstein’s relation between diffusion constant and conductivity, one gets the correct diffusion coefficient within the Boltzmann regime, $`D^{\text{(B)}}={\displaystyle \frac{3}{8}}lv`$.
To derive the results in Eq. (14), one uses Eq. (11) with $`\epsilon =0`$. The $`\epsilon `$-dependent term in Eq. (11) can be explicitly shown to be of higher order in the density than the terms present in Eq. (14) . In the following sections we will investigate the effect of sequences of correlated collisions between the point particle and the scatterers on the Lyapunov exponents. However, the $`\epsilon `$-dependent term in Eq. (11) will again be neglected since we will present the effect of these correlated collision sequences in leading order in the density of scatterers only. Thus, instead of Eq. (11), we will use
$`{\displaystyle \frac{1}{\rho _+}}={\displaystyle \frac{1}{\rho _{}}}+{\displaystyle \frac{2}{a\mathrm{cos}\varphi }}.`$ (16)
## 3 The extension of the ELBE to higher density
### 3.1 Binary collision operators in $`(\stackrel{}{r},\stackrel{}{v},\rho )`$-space and <br>the BBGKY hierarchy equations
The Boltzmann theory for the Lyapunov exponents assumes that the scatterers form a dilute, but quenched system and that the collisions of the point particles with the scatterers are uncorrelated. To incorporate the effects of correlated collisions on the Lyapunov exponents, we will use a method based on the BBGKY hierarchy equations, familiar from the kinetic theory of moderately dense gases . Since the moving particles do not interact with each other, it is sufficient to consider the distribution functions for just one of them, together with a number of scatterers. One starts from a fundamental equation for an $`(N+1)`$-body distribution function, $`f_{N+1}=f_{N+1}(\stackrel{}{r},\stackrel{}{v},\rho ;\stackrel{}{R}_1,\stackrel{}{R}_2,.,.,\stackrel{}{R}_N;t)`$, which is the probability density function in the entire extended phase space $`\mathrm{\Gamma }`$ spanned by the variables $`\stackrel{}{r},\stackrel{}{v},\rho ,\stackrel{}{R}_1,\stackrel{}{R}_2,.,.,\stackrel{}{R}_N`$, describing our system of $`N`$ scatterers and one moving particle. We require that $`f_{N+1}`$ satisfies the normalization condition
$`{\displaystyle }d\stackrel{}{r}d\stackrel{}{v}d\rho d\stackrel{}{R}_1d\stackrel{}{R}_2..d\stackrel{}{R}_Nf_{N+1}(\stackrel{}{r},\stackrel{}{v},\rho ;\stackrel{}{R}_1,\stackrel{}{R}_2,.,.,\stackrel{}{R}_N;t)=\mathrm{\hspace{0.17em}1}.`$ (17)
This $`(N+1)`$-body distribution function satisfies a Liouville-like equation determined by the collisions of the moving particles with the scatterers and by the motion of the particles in the thermostatted electric field, between collisions. Since the time evolution of $`\stackrel{}{r}`$ and $`\stackrel{}{v}`$ in this field is not Hamiltonian, we must use the Liouville equation in the form of a conservation law, rather than the usual form for Hamiltonian systems, to obtain
$`{\displaystyle \frac{f_{N+\mathrm{\hspace{0.17em}1}}}{t}}+\stackrel{}{}_\stackrel{}{r}(\dot{\stackrel{}{r}}f_{N+\mathrm{\hspace{0.17em}1}})+\stackrel{}{}_\stackrel{}{v}(\dot{\stackrel{}{v}}f_{N+\mathrm{\hspace{0.17em}1}})+{\displaystyle \frac{}{\rho }}(\dot{\rho }f_{N+\mathrm{\hspace{0.17em}1}})={\displaystyle \underset{i=\mathrm{\hspace{0.17em}1}}{\overset{N}{}}}\stackrel{~}{T}_{,i}f_{N+\mathrm{\hspace{0.17em}1}}.`$ (18)
Here the operators $`\stackrel{~}{T}_{,i}`$ are binary collision operators which describe the effects on the distribution function due to an instantaneous, elastic collision between the moving particle and the scatterer labeled by the index $`i`$. The explicit form of the binary collision operators may be easily obtained by a slight modification of the methods used by Ernst et al. , in order to include the radius of curvature as an additional variable. One finds that the action of this operator on any function $`f(\stackrel{}{r},\stackrel{}{v},\rho ;\stackrel{}{R}_1,\stackrel{}{R}_2,.,.,\stackrel{}{R}_j;t)`$ is
$`\stackrel{~}{T}_{,i}f=a{\displaystyle _{\stackrel{}{v}\widehat{\sigma }_i>\mathrm{\hspace{0.17em}0}}}d\widehat{\sigma }_i|\stackrel{}{v}\widehat{\sigma }_i|\{{\displaystyle _0^{\mathrm{}}}d\rho ^{}\delta (\rho {\displaystyle \frac{\rho ^{}a\mathrm{cos}\varphi _i}{a\mathrm{cos}\varphi _i+\mathrm{\hspace{0.17em}2}\rho ^{}}})\delta (a\widehat{\sigma }_i(\stackrel{}{r}\stackrel{}{R}_i))\times `$ (19)
$`\times b_{\sigma _i,\rho ^{}}\delta (a\widehat{\sigma }_i+(\stackrel{}{r}\stackrel{}{R}_i))\}f,`$
where $`\widehat{\sigma }_i`$ is the unit vector from the center of the scatterer fixed at $`\stackrel{}{R}_i`$ to the point of collision. The action of the operator $`b_{\sigma _i,\rho ^{}}`$ on the function $`f(\stackrel{}{r},\stackrel{}{v},\rho ;\stackrel{}{R}_1,\stackrel{}{R}_2,.,.,\stackrel{}{R}_j;t)`$ is defined by
$`b_{\sigma _i,\rho ^{}}f(\stackrel{}{r},\stackrel{}{v},\rho ;\stackrel{}{R}_1,\stackrel{}{R}_2,.,.,\stackrel{}{R}_j;t)=f(\stackrel{}{r},\stackrel{}{v}2(\stackrel{}{v}\widehat{\sigma }_i)\widehat{\sigma }_i,\rho ^{};\stackrel{}{R}_1,\stackrel{}{R}_2,.,.,\stackrel{}{R}_j;t);`$ (20)
that is, $`b_{\sigma _i\rho ^{}}`$ is a substitution operator that replaces $`\rho `$ by $`\rho ^{}`$ and the velocity $`\stackrel{}{v}`$ by its restituting value, i.e, the value it should have before collision so as to lead to the value $`\stackrel{}{v}`$ after collision. It is often useful to express the binary collision operators as as sum of two terms such that
$`\stackrel{~}{T}_{,i}=\stackrel{~}{T}_{,i}^{(+)}\stackrel{~}{T}_{,i}^{()},`$ (21)
where
$`\stackrel{~}{T}_{,i}^{(+)}=a{\displaystyle _{\stackrel{}{v}\widehat{\sigma }_i>\mathrm{\hspace{0.17em}0}}}𝑑\widehat{\sigma }_i|\stackrel{}{v}\widehat{\sigma }_i|{\displaystyle _0^{\mathrm{}}}𝑑\rho ^{}\delta \left(\rho {\displaystyle \frac{\rho ^{}a\mathrm{cos}\varphi _i}{a\mathrm{cos}\varphi _i+\mathrm{\hspace{0.17em}2}\rho ^{}}}\right)\delta (a\widehat{\sigma }_i(\stackrel{}{r}\stackrel{}{R}_i))b_{\sigma _i,\rho ^{}}`$ (22)
and
$`\stackrel{~}{T}_{,i}^{()}=a{\displaystyle _{\stackrel{}{v}\widehat{\sigma }_i>\mathrm{\hspace{0.17em}0}}}𝑑\widehat{\sigma }_i|\stackrel{}{v}\widehat{\sigma }_i|\delta (a\widehat{\sigma }_i+(\stackrel{}{r}\stackrel{}{R}_i)).`$ (23)
One sees that $`\stackrel{~}{T}_{,i}^{(+)}f`$ and $`\stackrel{~}{T}_{,i}^{()}f`$ respectively describe the rate of “gain” and the rate of “loss” of $`f`$ due to a collision of the point particle with the scatterer fixed at $`\stackrel{}{R}_i`$.
The BBGKY hierarchy equations are then obtained from Eq. (18) by integrating over scatterer coordinates, as a set of equations for the reduced distributions $`f_j`$ for the moving particle and $`(j1)`$ scatterers, defined by
$`f_j(\stackrel{}{r},\stackrel{}{v},\rho ;\stackrel{}{R}_1,\stackrel{}{R}_2,.,.,\stackrel{}{R}_{j1};t)`$ (24)
$`={\displaystyle \frac{N!}{(Nj+1)!}}{\displaystyle }d\stackrel{}{R}_j..d\stackrel{}{R}_Nf_{N+\mathrm{\hspace{0.17em}1}}(\stackrel{}{r},\stackrel{}{v},\rho ;\stackrel{}{R}_1,\stackrel{}{R}_2,.,.,\stackrel{}{R}_N;t).`$
One then easily obtains the BBGKY hierarchy equations ($`1jN`$)
$`{\displaystyle \frac{f_j}{t}}+\stackrel{}{}_\stackrel{}{r}(\dot{\stackrel{}{r}}f_j)+\stackrel{}{}_\stackrel{}{v}(\dot{\stackrel{}{v}}f_j)+{\displaystyle \frac{}{\rho }}(\dot{\rho }f_j){\displaystyle \underset{k=\mathrm{\hspace{0.17em}1}}{\overset{j\mathrm{\hspace{0.17em}1}}{}}}\stackrel{~}{T}_{,k}f_j={\displaystyle 𝑑\stackrel{}{R}_j\stackrel{~}{T}_{,j}f_{j+\mathrm{\hspace{0.17em}1}}}.`$ (25)
### 3.2 Cluster expansions and truncation of the hierarchy equations
The usual procedure for truncating the hierarchy equations in order to obtain the Boltzmann equation and its extension to higher densities is to make cluster expansions of the distribution functions, $`f_2,f_3...`$ in terms of a set of correlation functions, $`g_2,g_3...`$ as follows:
$`f_2(\stackrel{}{r},\stackrel{}{v},\rho ;\stackrel{}{R}_1;t)=nf_1(\stackrel{}{r},\stackrel{}{v},\rho ;t)+g_2(\stackrel{}{r},\stackrel{}{v},\rho ;\stackrel{}{R}_1;t),`$ (26)
$`f_3(\stackrel{}{r},\stackrel{}{v},\rho ;\stackrel{}{R}_1,\stackrel{}{R}_2;t)=n^2f_1(\stackrel{}{r},\stackrel{}{v},\rho ;t)+ng_2(\stackrel{}{r},\stackrel{}{v},\rho ;\stackrel{}{R}_1;t)+ng_2(\stackrel{}{r},\stackrel{}{v},\rho ;\stackrel{}{R}_2;t)`$ (27)
$`+g_3(\stackrel{}{r},\stackrel{}{v},\rho ;\stackrel{}{R}_1,\stackrel{}{R}_2;t),`$
and so on. Hereafter, to save writing, we denote $`g_2(\stackrel{}{r},\stackrel{}{v},\rho ;\stackrel{}{R}_1;t)`$ as $`g_{2,\stackrel{}{R}_1}`$, $`g_2(\stackrel{}{r},\stackrel{}{v},\rho ;\stackrel{}{R}_2;t)`$ as $`g_{2,\stackrel{}{R}_2}`$, $`f_3(\stackrel{}{r},\stackrel{}{v},\rho ;\stackrel{}{R}_1,\stackrel{}{R}_2;t)`$ as $`f_3`$ and $`g_3(\stackrel{}{r},\stackrel{}{v},\rho ;\stackrel{}{R}_1,\stackrel{}{R}_2;t)`$ as $`g_3`$. The first terms in each of these expansions represent the totally uncorrelated situation, where there are independent probabilities of finding the moving particle and the scatterers at the designated coordinates. The next terms involving the pair correlation functions $`g_{2,\stackrel{}{R}_i}`$ in Eqs. (26) and (27) take into account possible dynamical and excluded volume correlations between the point particle and the scatterer at $`\stackrel{}{R}_i`$. If one replaces $`f_2`$ by $`nf_1`$ in the first BBGKY hierarchy equation, Eq. (25) with $`j=1`$, reduces to the ELBE. To find the corrections to the ELBE for higher densities, one must keep the $`g_2`$ term in Eq. (26) and use the second hierarchy equation to determine $`g_2`$. However, in order to solve the second equation, we have to say something about $`g_3`$. A careful examination of the second and higher equations shows that $`g_3`$ contains, of course, the effects of three-body correlations, i.e, correlated collisions involving the point particle, a scatterer fixed at $`\stackrel{}{R}_1`$ and another scatterer fixed at $`\stackrel{}{R}_2`$, as well as excluded volume corrections due to the non-overlapping property of the scatterers. Here we will be primarily interested in the effects of the so called “ring” collisions on the Lyapunov exponents. These collision sequences are composed of one collision of the moving particle with a given scatterer, followed by an arbitrary number of collisions with a succession of different scatterers, and completed by a final re-collision of the moving particle with the first scatterer in the sequence, as illustrated in Fig 3.
Fig. 3 : Sequential collisions with scatterers at $`\stackrel{}{R}_2,\stackrel{}{R}_3,.,.`$ adding up to the ring diagram.
The ring diagrams, taken individually, are the most divergent terms that appear in the expansion of dynamical properties of the Lorentz gas as a power series in the density of scatterers. They lead to the logarithmic terms in the density expansion of the diffusion coefficient of the moving particle , and to the algebraic long time tails in the velocity time correlation function of the moving particle . While many other dynamical events and excluded volume effects contribute to the Lyapunov exponents, and must be included for a full treatment, we concentrate here on the effects of these most divergent collision sequences, since in other contexts, they are responsible for the most dramatic higher density corrections to the Boltzmann equation results.
Thus we drop $`g_3`$ in Eq. (27) and obtain a somewhat simplified cluster expansion of $`f_3`$, given by
$`f_3=n^2f_1+ng_{2,\stackrel{}{R}_1}+ng_{2,\stackrel{}{R}_2}`$ (28)
Using Eqs. (26) and (28) and the first two of the BBGKY hierarchy equations, we obtain a closed set of two equations involving two unknowns, $`f_1`$ and $`g_{2,\stackrel{}{R}_1}`$ given by
$`\stackrel{}{}_\stackrel{}{v}(\dot{\stackrel{}{v}}f_1)+{\displaystyle \frac{}{\rho }}(\dot{\rho }f_1)={\displaystyle 𝑑\stackrel{}{R}_1\stackrel{~}{T}_{,\mathrm{\hspace{0.17em}1}}[nf_1+g_{2,\stackrel{}{R}_1}]}`$ (29)
and
$`\stackrel{}{}_\stackrel{}{r}(\dot{\stackrel{}{r}}g_{2,\stackrel{}{R}_1})+\stackrel{}{}_\stackrel{}{v}(\dot{\stackrel{}{v}}g_{2,\stackrel{}{R}_1})+{\displaystyle \frac{}{\rho }}(\dot{\rho }g_{2,\stackrel{}{R}_1})n{\displaystyle 𝑑\stackrel{}{R}_2\stackrel{~}{T}_{,\mathrm{\hspace{0.17em}2}}g_{2,\stackrel{}{R}_1}}=n\stackrel{~}{T}_{,\mathrm{\hspace{0.17em}1}}f_1.`$ (30)
In the derivation of Eq. (30) from the second hierarchy equation not only have we dropped $`g_3`$ as discussed above, we also dropped a term of the form $`\stackrel{~}{T}_{,\mathrm{\hspace{0.17em}1}}g_{2,\stackrel{}{R_1}}`$. This term provides “repeated ring” corrections to the ring contributions to $`g_{2,\stackrel{}{R_1}}`$. These are of the same order as terms neglected by dropping $`g_3`$ , and should be neglected for consistency. We also dropped the time derivatives in the equations, so we are now looking for the distribution and correlation functions appropriate for the NESS.
In Section 4, we will solve Eqs. (29) and (30) in order to calculate the ring contributions to the positive Lyapunov exponent. Before doing so, it is useful to write down the usual form of the ring equations in $`(\stackrel{}{r},\stackrel{}{v})`$-space, which can be obtained by integrating Eqs. (29) and (30) over all values of the radius of curvature, $`0\rho <\mathrm{}`$. We define the usual single-particle distribution function by, $`F_1=_{\rho >0}𝑑\rho f_1`$ and the pair-correlation function $`G_{2,\stackrel{}{R}_1}=_{\rho >0}𝑑\rho g_{2,\stackrel{}{R}_1}`$. By imposing the boundary conditions that both $`f_1`$ and $`g_{2,\stackrel{}{R}_1}`$ go to zero as $`\rho 0`$ and as $`\rho \mathrm{}`$, we obtain
$`\stackrel{}{}_\stackrel{}{v}(\dot{\stackrel{}{v}}F_1)={\displaystyle 𝑑\stackrel{}{R}_1\overline{T}_{,\mathrm{\hspace{0.17em}1}}[nF_1+G_{2,\stackrel{}{R}_1}]}`$ (31)
and
$`\stackrel{}{}_\stackrel{}{r}(\dot{\stackrel{}{r}}G_{2,\stackrel{}{R}_1})+\stackrel{}{}_\stackrel{}{v}(\dot{\stackrel{}{v}}G_{2,\stackrel{}{R}_1})n{\displaystyle 𝑑\stackrel{}{R}_2\overline{T}_{,\mathrm{\hspace{0.17em}2}}G_{2,\stackrel{}{R}_1}}=n\overline{T}_{,\mathrm{\hspace{0.17em}1}}F_1.`$ (32)
The actions of $`\overline{T}_{,\mathrm{\hspace{0.17em}1}}`$ or $`\overline{T}_{,\mathrm{\hspace{0.17em}2}}`$ on $`F_1`$ and $`G_{2,\stackrel{}{R}_1}`$ can be obtained by appropriately integrating $`\stackrel{~}{T}_{,\mathrm{\hspace{0.17em}1}}f_1`$, $`\stackrel{~}{T}_{,\mathrm{\hspace{0.17em}1}}g_{2,\stackrel{}{R}_1}`$ or $`\stackrel{~}{T}_{,\mathrm{\hspace{0.17em}2}}g_{2,\stackrel{}{R}_1}`$ over $`\rho `$ from $`0`$ to $`\mathrm{}`$ using the definitions in Eqs. (19) and (20). $`\overline{T}_{,\mathrm{\hspace{0.17em}1}}`$ and $`\overline{T}_{,\mathrm{\hspace{0.17em}2}}`$ are the analogs in ($`\stackrel{}{r},\stackrel{}{v}`$) space of $`\stackrel{~}{T}_{,\mathrm{\hspace{0.17em}1}}`$ and $`\stackrel{~}{T}_{,\mathrm{\hspace{0.17em}2}}`$ (see Eqs. (19) and (20)), i.e.,
$`\overline{T}_{,i}=a{\displaystyle _{\stackrel{}{v}\widehat{\sigma }_i>\mathrm{\hspace{0.17em}0}}}𝑑\widehat{\sigma }_i|\stackrel{}{v}\widehat{\sigma }_i|\left\{\delta (a\widehat{\sigma }_i(\stackrel{}{r}\stackrel{}{R}_i))b_{\sigma _i}\delta (a\widehat{\sigma }_i+(\stackrel{}{r}\stackrel{}{R}_i))\right\}.`$ (33)
In future applications however, we will drop the $`a\widehat{\sigma }_i`$ terms from the arguments of both $`\delta (a\widehat{\sigma }_i\pm (\stackrel{}{r}\stackrel{}{R}_i))`$ in $`\stackrel{~}{T}_{,i}`$ and $`\overline{T}_{,i}`$ operators since they lead to corrections similar to excluded volume terms, neglected already.
## 4 Effects of long range time correlation on $`\lambda _+`$ and $`\lambda _{}`$
We now concentrate on the solution of the BBGKY equations for the distribution functions that determine the Lyapunov exponents. The solutions of Eqs. (29) and (30) are to be obtained as expansions in two small variables, $`na^2`$ and $`\stackrel{~}{\epsilon }`$. The density expansion will give the corrections to the previously obtained Boltzmann regime results from the ELBE, and the $`\stackrel{~}{\epsilon }`$ expansion will provide the field dependence of these corrections. We therefore write the density expansions of $`f_1`$ and $`g_2`$ (hereafter we drop the subscript $`\stackrel{}{R}_1`$ from $`g_{2,\stackrel{}{R}_1}`$) as
$`f_1=f_1^{\text{(B)}}+f_1^{\text{(R)}}+...\text{and}g_2=g_2^{\text{(R)}}+...,`$ (34)
where the superscript B indicates the lowest density result for $`f_1`$ as given by the ELBE, and the superscript R denotes the ring contribution. At the order in density of interest here, the quantities indicated explicitly in the above equations satisfy
$`\stackrel{}{}_\stackrel{}{v}(\dot{\stackrel{}{v}}f_1^{\text{(B)}})+{\displaystyle \frac{}{\rho }}(\dot{\rho }f_1^{\text{(B)}})=n{\displaystyle 𝑑\stackrel{}{R}_1\stackrel{~}{T}_{,\mathrm{\hspace{0.17em}1}}f_1^{\text{(B)}}},`$ (35)
$`\stackrel{}{}_\stackrel{}{v}(\dot{\stackrel{}{v}}f_1^{\text{(R)}})+{\displaystyle \frac{}{\rho }}(\dot{\rho }f_1^{\text{(R)}})={\displaystyle 𝑑\stackrel{}{R}_1\stackrel{~}{T}_{,\mathrm{\hspace{0.17em}1}}[nf_1^{\text{(R)}}+g_2^{\text{(R)}}]}`$ (36)
and
$`\stackrel{}{}_\stackrel{}{r}(\dot{\stackrel{}{r}}g_2^{\text{(R)}})+\stackrel{}{}_\stackrel{}{v}(\dot{\stackrel{}{v}}g_2^{\text{(R)}})+{\displaystyle \frac{}{\rho }}(\dot{\rho }g_2^{\text{(R)}})n{\displaystyle 𝑑\stackrel{}{R}_2\stackrel{~}{T}_{,\mathrm{\hspace{0.17em}2}}g_2^{\text{(R)}}}=n\stackrel{~}{T}_{,\mathrm{\hspace{0.17em}1}}f_1^{\text{(B)}}.`$ (37)
Our aim here is to solve Eqs. (36) and (37) using the results of the ELBE for $`f_1^{\text{(B)}}`$. We suppose further that each of these functions possesses an expansion in powers of $`\epsilon `$ as
$`f_1^{\text{(B,\hspace{0.17em}R)}}`$ $`=`$ $`f_{1,\mathrm{\hspace{0.17em}0}}^{\text{(B,\hspace{0.17em}R)}}+\epsilon f_{1,\mathrm{\hspace{0.17em}1}}^{\text{(B,\hspace{0.17em}R)}}+\epsilon ^2f_{1,\mathrm{\hspace{0.17em}2}}^{\text{(B,\hspace{0.17em}R)}}+...`$ (38)
and
$`g_2^{\text{(R)}}`$ $`=`$ $`g_{2,\mathrm{\hspace{0.17em}0}}^{\text{(R)}}+\epsilon g_{2,\mathrm{\hspace{0.17em}1}}^{\text{(R)}}+\epsilon ^2g_{2,\mathrm{\hspace{0.17em}2}}^{\text{(R)}}+...`$ (39)
The functions $`f_{1,i}^{\text{(B)}}`$ have been previously obtained as the $`\epsilon `$ solutions of the ELBE. Since we will be dealing with $`g_2`$ only in the context of of the ring term, we drop the superscript R from now on.
As mentioned above, we will neglect the term $`a\widehat{\sigma }`$ within the arguments of the $`\delta `$-functions appearing in each of the binary collision operators $`\stackrel{~}{T}_{}`$ and $`\overline{T}_{}`$, so as to take the moving particle to be located at the same point as the center of the appropriate scatterer at collision. The terms neglected by this approximation lead to higher density corrections to the terms we will obtain below. Secondly, an inspection of the radius of curvature delta function in the expression for the “gain” part of the binary collision operator, Eq. (22), shows that this term is only non-vanishing when $`\rho {\displaystyle \frac{a}{2}}`$, and that the dominant contribution to the $`\rho ^{}`$ integration comes from the region $`\rho ^{}l`$. Naturally, $`{\displaystyle \frac{\rho ^{}a\mathrm{cos}\varphi _i}{a\mathrm{cos}\varphi _i+\mathrm{\hspace{0.17em}2}\rho ^{}}}{\displaystyle \frac{a\mathrm{cos}\varphi _i}{2}}(1+O(n))`$ in the argument of the delta function. In the Boltzmann level approximation this $`O(n)`$ term may therefore be neglected and it can be shown not to contribute to the leading field-dependent ring term effects on the Lyapunov exponents. Therefore, we will neglect it in what follows. Under this approximation, the gain part of the binary collision operator acts on an arbitrary function $`h(\stackrel{}{r},\stackrel{}{v},\rho )`$ as
$`\stackrel{~}{T}_{,i}^{(+)}h(\stackrel{}{r},\stackrel{}{v},\rho )\delta (\stackrel{}{r}\stackrel{}{R}_i)\mathrm{\Theta }\left({\displaystyle \frac{a}{2}}\rho \right)I(\rho )[H(\stackrel{}{r},\stackrel{}{v}_+)+H(\stackrel{}{r},\stackrel{}{v}_{})]\delta (\stackrel{}{r}\stackrel{}{R}_i)\mathrm{\Gamma }(\rho ,H).`$ (40)
Here
$`\stackrel{}{v}_\pm =\stackrel{}{v}\mathrm{\hspace{0.17em}2}(\stackrel{}{v}\widehat{\sigma }_{i,\pm })\widehat{\sigma }_{i,\pm },`$ (41)
and $`\widehat{\sigma }_{i,\pm }`$ is defined by the condition that the scattering angle $`\varphi =\pm \mathrm{cos}^1\left({\displaystyle \frac{2\rho }{a}}\right)`$. Also
$`H(\stackrel{}{r},\stackrel{}{v})={\displaystyle _0^{\mathrm{}}}𝑑\rho ^{}h(\stackrel{}{r},\stackrel{}{v},\rho ^{})`$ (42)
and
$$I(\rho )=\frac{4v\rho }{a\sqrt{1\left(\frac{2\rho }{a}\right)^2}}.$$
(43)
Finally, we express the velocity vector $`\stackrel{}{v}`$ in terms of the angle $`\theta `$ that it makes with the direction of the electric field so as to obtain the following set of equations for the terms in the $`\epsilon `$-expansion of $`g_2`$
$`\left[\stackrel{}{v}\stackrel{}{}_\stackrel{}{r}+v{\displaystyle \frac{}{\rho }}+\mathrm{\hspace{0.17em}2}nav\right]g_{2,\mathrm{\hspace{0.17em}0}}=\mathrm{\hspace{0.17em}2}nav\delta (\stackrel{}{r}\stackrel{}{R}_1)f_{1,\mathrm{\hspace{0.17em}0}}^{\text{(B)}}+n\delta (\stackrel{}{r}\stackrel{}{R}_1)\mathrm{\Gamma }(\rho ,F_{1,0}^{\text{(B)}}),`$ (44)
$`\left[\stackrel{}{v}\stackrel{}{}_\stackrel{}{r}+v{\displaystyle \frac{}{\rho }}+\mathrm{\hspace{0.17em}2}nav\right]g_{2,\mathrm{\hspace{0.17em}1}}={\displaystyle \frac{}{\rho }}\left(\rho \mathrm{cos}\theta g_{2,\mathrm{\hspace{0.17em}0}}\right)+{\displaystyle \frac{}{\theta }}\left(\mathrm{sin}\theta g_{2,\mathrm{\hspace{0.17em}0}}\right)`$ (45)
$`\mathrm{\hspace{0.17em}2}nav\delta (\stackrel{}{r}\stackrel{}{R}_1)f_{1,\mathrm{\hspace{0.17em}1}}^{\text{(B)}}+n\delta (\stackrel{}{r}\stackrel{}{R}_1)\mathrm{\Gamma }(\rho ,F_{1,1}^{\text{(B)}}),`$
and
$`\left[\stackrel{}{v}\stackrel{}{}_\stackrel{}{r}+v{\displaystyle \frac{}{\rho }}+\mathrm{\hspace{0.17em}2}nav\right]g_{2,\mathrm{\hspace{0.17em}2}}={\displaystyle \frac{}{\rho }}\left(\rho \mathrm{cos}\theta g_{2,\mathrm{\hspace{0.17em}1}}+{\displaystyle \frac{\rho ^2\mathrm{sin}^2\theta }{v}}g_{2,\mathrm{\hspace{0.17em}0}}\right)+{\displaystyle \frac{}{\theta }}\left(\mathrm{sin}\theta g_{2,\mathrm{\hspace{0.17em}1}}\right)`$ (46)
$`\mathrm{\hspace{0.17em}2}nav\delta (\stackrel{}{r}\stackrel{}{R}_1)f_{1,\mathrm{\hspace{0.17em}2}}^{\text{(B)}}+n\delta (\stackrel{}{r}\stackrel{}{R}_1)\mathrm{\Gamma }(\rho ,F_{1,2}^{\text{(B)}}).`$
We notice that the equations thus obtained are linear inhomogeneous differential equations of the form $`Lg_{2,j}=b_j`$ ($`j=0,1,2`$), where $`L=\left[\stackrel{}{v}\stackrel{}{}_\stackrel{}{r}+v{\displaystyle \frac{}{\rho }}+\mathrm{\hspace{0.17em}2}nav\right]`$ is a linear differential operator and $`b_j`$ for $`j=0,1,2`$ are the inhomogeneous terms on the r.h.s. of Eqs. (44), (45) and(46) respectively.
We will need to solve Eqs. (44-46), in conjunction with the equations for $`G_{2,\mathrm{\hspace{0.17em}0}}`$, $`G_{2,\mathrm{\hspace{0.17em}1}}`$ and $`G_{2,\mathrm{\hspace{0.17em}2}}`$ obtained by directly integrating Eqs. (44-46) over $`\rho `$ from 0 to $`\mathrm{}`$. The equations for the corresponding $`G_{2,i}`$ are then
$`\left[\stackrel{}{v}\stackrel{}{}_\stackrel{}{r}n{\displaystyle \stackrel{}{dR_2}\overline{T}_{,\mathrm{\hspace{0.17em}2}}}\right]G_{2,\mathrm{\hspace{0.17em}0}}`$ $`=`$ $`n\overline{T}_{,\mathrm{\hspace{0.17em}1}}F_{1,\mathrm{\hspace{0.17em}0}}^{\text{(B)}},`$ (47)
$`\left[\stackrel{}{v}\stackrel{}{}_\stackrel{}{r}n{\displaystyle \stackrel{}{dR_2}\overline{T}_{,\mathrm{\hspace{0.17em}2}}}\right]G_{2,\mathrm{\hspace{0.17em}1}}`$ $`=`$ $`{\displaystyle \frac{}{\theta }}\left(\mathrm{sin}\theta G_{2,\mathrm{\hspace{0.17em}0}}\right)+n\overline{T}_{,\mathrm{\hspace{0.17em}1}}F_{1,\mathrm{\hspace{0.17em}1}}^{\text{(B)}}`$ (48)
and
$`\left[\stackrel{}{v}\stackrel{}{}_\stackrel{}{r}n{\displaystyle \stackrel{}{dR_2}\overline{T}_{,\mathrm{\hspace{0.17em}2}}}\right]G_{2,\mathrm{\hspace{0.17em}2}}`$ $`=`$ $`{\displaystyle \frac{}{\theta }}\left(\mathrm{sin}\theta G_{2,\mathrm{\hspace{0.17em}1}}\right)+n\overline{T}_{,\mathrm{\hspace{0.17em}1}}F_{1,\mathrm{\hspace{0.17em}2}}^{\text{(B)}}.`$ (49)
The equations for $`G_{2,\mathrm{\hspace{0.17em}0}}`$, $`G_{2,\mathrm{\hspace{0.17em}1}}`$ and $`G_{2,\mathrm{\hspace{0.17em}2}}`$ are also linear differential equations of the form $`L^{\frac{}{}}G_{2,j}=B_j`$, where $`L^{\frac{}{}}=\left[\stackrel{}{v}\stackrel{}{}_\stackrel{}{r}n\stackrel{}{dR_2}\overline{T}_{,\mathrm{\hspace{0.17em}2}}\right]`$ and $`B_j`$ for $`j=\mathrm{\hspace{0.17em}0},1,2`$ are the inhomogeneous terms on the r.h.s. of Eqs. (47), (48) and (49) respectively.
To solve these equations, we take Fourier transforms of $`g_2`$ and of $`G_2`$ in the variables $`\stackrel{}{r}`$ and in the velocity angle, $`\theta `$. That is, we define $`\stackrel{~}{g}_2(\stackrel{}{k})={\displaystyle \frac{1}{\sqrt{V}}}{\displaystyle _V}\stackrel{}{dr}g_2e^{i\stackrel{}{k}(\stackrel{}{r}\stackrel{}{R}_1)}`$ and calculate $`\stackrel{~}{g}_{2,\mathrm{\hspace{0.17em}0}}(\stackrel{}{k})`$, $`\stackrel{~}{g}_{2,\mathrm{\hspace{0.17em}1}}(\stackrel{}{k})`$ and $`\stackrel{~}{g}_{2,\mathrm{\hspace{0.17em}2}}(\stackrel{}{k})`$ in the $`\stackrel{}{k}`$-basis, using periodic boundary conditions. Similarly, we define the $`m`$-th angular mode of $`\stackrel{~}{g}_2(\stackrel{}{k})`$ as $`\stackrel{~}{g}_2^{(m)}(\stackrel{}{k})={\displaystyle \frac{1}{\sqrt{2\pi }}}{\displaystyle 𝑑\theta e^{im\theta }\stackrel{~}{g}_2(\stackrel{}{k})}`$. We also define $`\stackrel{~}{G}_2^{(m)}(\stackrel{}{k})`$ in an analogous way. Thus, corresponding to Eqs. (44-46) and Eqs. (47-49), we have two sets of three equations to be solved, one involving $`\stackrel{~}{g}_{2,\mathrm{\hspace{0.17em}0}}(\stackrel{}{k})`$, $`\stackrel{~}{g}_{2,\mathrm{\hspace{0.17em}1}}(\stackrel{}{k})`$ and $`\stackrel{~}{g}_{2,\mathrm{\hspace{0.17em}2}}(\stackrel{}{k})`$, and the other involving $`\stackrel{~}{G}_{2,\mathrm{\hspace{0.17em}0}}(\stackrel{}{k})`$, $`\stackrel{~}{G}_{2,\mathrm{\hspace{0.17em}1}}(\stackrel{}{k})`$ and $`\stackrel{~}{G}_{2,\mathrm{\hspace{0.17em}2}}(\stackrel{}{k})`$, in $`(\stackrel{}{k},m)`$ basis. In this basis, the operators $`L_\stackrel{}{k}=\left[i\stackrel{}{k}\stackrel{}{v}+v{\displaystyle \frac{}{\rho }}+\mathrm{\hspace{0.17em}2}nav\right]`$ and $`L_\stackrel{}{k}^{\frac{}{}}=\left[i\stackrel{}{k}\stackrel{}{v}n\stackrel{}{dR_2}\overline{T}_{,\mathrm{\hspace{0.17em}2}}\right]`$ are both infinite dimensional matrices in $`m`$-space and both of them have non-zero off-diagonal elements due to the term $`i\stackrel{}{k}\stackrel{}{v}`$ generated from the operator $`\stackrel{}{v}\stackrel{}{}_\stackrel{}{r}`$. However, it is easily seen that these off-diagonal elements are proportional to $`\delta _{m,m+1}`$ and $`\delta _{m,m1}`$ and they are easily treated.
A further simplification can be made by noticing that the schematic forms of the solutions are $`\stackrel{~}{g}_{2,j}(\stackrel{}{k})=[L_\stackrel{}{k}]^1b_j(\stackrel{}{k})`$ and $`\stackrel{~}{G}_{2,j}(\stackrel{}{k})=[L_\stackrel{}{k}^{\frac{}{}}]^1B_j(\stackrel{}{k})`$ and hence the dominant parts of $`\stackrel{~}{g}_{2,j}(\stackrel{}{k})`$ and $`\stackrel{~}{G}_{2,j}(\stackrel{}{k})`$ will come, loosely speaking, from the eigenfunctions of $`L_\stackrel{}{k}`$ and $`L_\stackrel{}{k}^{\frac{}{}}`$ having the smallest eigenvalues. The lowest eigenvalues of $`L_\stackrel{}{k}^{\frac{}{}}`$ are $`k^2`$ due to the contributions from the hydrodynamic modes . Thus, to capture the dominant part of the solutions we should solve the equations in the range $`k=|\stackrel{}{k}|<<l^1`$, the inverse mean free path and use perturbation expansions in the small parameter $`kl`$. We will not, in our analysis, follow the mode expansion technique, as it is simpler to calculate $`G_2`$ directly. However, one can use mode expansions and one finds that the results of both the methods agree.
### 4.1 Solution for $`G_2`$
To solve for $`G_2`$ first we need to know the solutions of the Lorentz-Boltzmann equation for $`F_{1,\mathrm{\hspace{0.17em}0}}^{\text{(B)}}`$, $`F_{1,\mathrm{\hspace{0.17em}1}}^{\text{(B)}}`$ and $`F_{1,\mathrm{\hspace{0.17em}2}}^{\text{(B)}}`$. These are given by
$`F_{1,\mathrm{\hspace{0.17em}0}}^{\text{(B)}}={\displaystyle \frac{1}{2\pi }},F_{1,\mathrm{\hspace{0.17em}1}}^{\text{(B)}}={\displaystyle \frac{3}{16\pi nav}}\mathrm{cos}\theta \text{and}F_{1,\mathrm{\hspace{0.17em}2}}^{\text{(B)}}={\displaystyle \frac{45}{512\pi (nav)^2}}\mathrm{cos}2\theta .`$ (50)
We note that in $`m`$-space, defined above, the $`m`$-th diagonal element of the infinite matrix $`L_\stackrel{}{k}^{^{}}`$ is $`{\displaystyle \frac{4m^2}{(4m^2\mathrm{\hspace{0.17em}1})}}{\displaystyle \frac{v}{l}}`$ while the off-diagonal elements are $`ikv\delta _{m,m\pm 1}`$. Thus an expansion in $`\stackrel{~}{k}=kl`$ can be easily obtained by considering successively larger parts of the matrix $`L_\stackrel{}{k}^{^{}}`$ in the index $`m`$, starting with $`3\times 3`$, $`5\times 5`$ matrices and so on, chosen in such a way that the element of $`L_\stackrel{}{k}^{^{}}`$ corresponding to $`m=0`$ appears as the center element of these matrices. As we want to make our results correct up to $`O(\stackrel{~}{k}^0)`$, we need to increase the size of these matrices till the expressions of $`\stackrel{~}{G}_{2,j}(\stackrel{}{k})`$ obtained from $`\stackrel{~}{G}_{2,j}(\stackrel{}{k})=[L_\stackrel{}{k}^{\frac{}{}}]^1B_j(\stackrel{}{k})`$ (for $`j=0,1,2`$) converges up to $`O(\stackrel{~}{k}^0)`$. Also, as we want to obtain the expression of $`\lambda _+`$ and $`\lambda _{}`$ in the leading field-dependent order, which is $`\epsilon ^2`$; we need the solutions of all the $`m`$-modes of $`\stackrel{~}{G}_{2,\mathrm{\hspace{0.17em}0}}(\stackrel{}{k})`$, $`\stackrel{~}{G}_{2,\mathrm{\hspace{0.17em}1}}(\stackrel{}{k})`$ and $`\stackrel{~}{G}_{2,\mathrm{\hspace{0.17em}2}}(\stackrel{}{k})`$ that are necessary to obtain all the terms of $`f_1^{\text{(R)}}`$ that are $`\epsilon ^2`$ and contribute to this leading field-dependent order of $`\lambda _+`$ and $`\lambda _{}`$. In more explicit form, this means that we definitely need the solutions of $`\stackrel{~}{G}_{2,\mathrm{\hspace{0.17em}0}}^{(m=\mathrm{\hspace{0.17em}0})}(\stackrel{}{k})`$, $`\stackrel{~}{G}_{2,\mathrm{\hspace{0.17em}1}}^{(m=\mathrm{\hspace{0.17em}0})}(\stackrel{}{k})`$, $`\stackrel{~}{G}_{2,\mathrm{\hspace{0.17em}1}}^{(m=\pm 1)}(\stackrel{}{k})`$, $`\stackrel{~}{G}_{2,\mathrm{\hspace{0.17em}1}}^{(m=\pm 2)}(\stackrel{}{k})`$ and $`\stackrel{~}{G}_{2,\mathrm{\hspace{0.17em}2}}^{(m=\mathrm{\hspace{0.17em}0})}(\stackrel{}{k})`$ up to $`O(\stackrel{~}{k}^0)`$. However, once we present these solutions, from the structure and properties of them, it will turn out that we will also need the expressions of $`\stackrel{~}{G}_{2,\mathrm{\hspace{0.17em}2}j}^{(m=\mathrm{\hspace{0.17em}0})}(\stackrel{}{k})`$ in the leading order of $`\stackrel{~}{k}`$ for $`j=\mathrm{\hspace{0.17em}2},3,..`$, to consistently obtain all the terms, that are $`\epsilon ^2`$.
At the $`\epsilon ^0`$ or equilibrium order, we find
$`\stackrel{~}{G}_{2,\mathrm{\hspace{0.17em}0}}^{(m)}(\stackrel{}{k})=\mathrm{\hspace{0.17em}0}m.`$ (51)
Proceeding to order $`\epsilon `$, we find that $`\stackrel{~}{G}_{2,\mathrm{\hspace{0.17em}1}}^{(m)}(\stackrel{}{k})`$’s obey
$`{\displaystyle \frac{iv}{2}}\left[(k_x+ik_y)\stackrel{~}{G}_{2,\mathrm{\hspace{0.17em}1}}^{(m+1)}(\stackrel{}{k})+(k_xik_y)\stackrel{~}{G}_{2,\mathrm{\hspace{0.17em}1}}^{(m1)}(\stackrel{}{k})\right]+{\displaystyle \frac{8navm^2}{4m^2\mathrm{\hspace{0.17em}1}}}\stackrel{~}{G}_{2,\mathrm{\hspace{0.17em}1}}^{(m)}(\stackrel{}{k})`$ (52)
$`={\displaystyle \frac{1}{2\sqrt{2\pi V}}}\left(\delta _{m,\mathrm{\hspace{0.17em}1}}+\delta _{m,1}\right),`$
where it turns out that we need a $`5\times 5`$ matrix block corresponding to $`m=2,1,\mathrm{\hspace{0.17em}0},\mathrm{\hspace{0.17em}1},`$ and $`2`$ to get the solutions of $`\stackrel{~}{G}_{2,\mathrm{\hspace{0.17em}1}}^{(m)}(\stackrel{}{k})`$ up to $`O(k^0)`$ that are relevant for us, yielding
$`\stackrel{~}{G}_{2,\mathrm{\hspace{0.17em}1}}^{(m=\mathrm{\hspace{0.17em}0})}(\stackrel{}{k})={\displaystyle \frac{ik_x}{vk^2\sqrt{2\pi V}}},`$
$`\stackrel{~}{G}_{2,\mathrm{\hspace{0.17em}1}}^{(m=1)}(\stackrel{}{k})={\displaystyle \frac{3ik_y(k_xik_y)}{16navk^2\sqrt{2\pi V}}},\stackrel{~}{G}_{2,\mathrm{\hspace{0.17em}1}}^{(m=1)}(\stackrel{}{k})={\displaystyle \frac{3ik_y(k_x+ik_y)}{16navk^2\sqrt{2\pi V}}}`$
$`\stackrel{~}{G}_{2,\mathrm{\hspace{0.17em}1}}^{(m=\mathrm{\hspace{0.17em}2})}(\stackrel{}{k})={\displaystyle \frac{45k_y(k_xik_y)^2}{1024(na)^2vk^2\sqrt{2\pi V}}}\text{and}\stackrel{~}{G}_{2,\mathrm{\hspace{0.17em}1}}^{(m=2)}(\stackrel{}{k})={\displaystyle \frac{45k_y(k_x+ik_y)^2}{1024(na)^2vk^2\sqrt{2\pi V}}}.`$ (53)
Notice that we have also calculated $`\stackrel{~}{G}_{2,\mathrm{\hspace{0.17em}1}}^{(m=\mathrm{\hspace{0.17em}2})}(\stackrel{}{k})`$ and $`\stackrel{~}{G}_{2,\mathrm{\hspace{0.17em}1}}^{(m=\mathrm{\hspace{0.17em}2})}(\stackrel{}{k})`$, even though they are $`O(k)`$, because they affect the $`O(k^0)`$ solution for $`\stackrel{~}{G}_{2,\mathrm{\hspace{0.17em}2}}^{(m=\mathrm{\hspace{0.17em}0})}(\stackrel{}{k})`$.
For order $`\epsilon ^2`$, the relevant $`\stackrel{~}{G}_{2,\mathrm{\hspace{0.17em}2}}^{(m)}(\stackrel{}{k})`$’s are then calculated using Eq. (53) and considering a $`5\times 5`$ matrix block of $`L_\stackrel{}{k}^{\frac{}{}}`$. There we need only the solution for $`\stackrel{~}{G}_{2,\mathrm{\hspace{0.17em}2}}^{(m=\mathrm{\hspace{0.17em}0})}(\stackrel{}{k})`$ :
$`\stackrel{~}{G}_{2,\mathrm{\hspace{0.17em}2}}^{(m=\mathrm{\hspace{0.17em}0})}(\stackrel{}{k})={\displaystyle \frac{k_x^2}{v^2k^4\sqrt{2\pi V}}}+{\displaystyle \frac{45(2k_x^2\mathrm{\hspace{0.17em}5}k_y^2)}{1024(nav)^2k^2\sqrt{2\pi V}}}.`$ (54)
Examining the properties of the solutions, Eqs. (53-54) and observing from Eqs. (47-49) the way the solution of $`G_{2,j}`$ affects the solution of $`G_{2,(j+\mathrm{\hspace{0.17em}1})}`$, one sees that the leading power of $`k`$ in the expression of $`\stackrel{~}{G}_{2,\mathrm{\hspace{0.17em}2}j}^{(m=\mathrm{\hspace{0.17em}0})}(\stackrel{}{k})`$ for ($`j=1,2,3\mathrm{}.`$) is $`k^{\mathrm{\hspace{0.17em}2}j}`$. However, in the expression of $`G_2`$, $`\stackrel{~}{G}_{2,\mathrm{\hspace{0.17em}2}j}^{(m=\mathrm{\hspace{0.17em}0})}(\stackrel{}{k})`$ appears with a factor of $`\epsilon ^{2j}`$. When $`G_2`$ is finally calculated, after a summation of the appropriate $`\stackrel{}{k}`$-values<sup>1</sup><sup>1</sup>1To see how the $`\stackrel{}{k}`$-integration is performed, see the last paragraph of Section 4.2, the contribution of the sum of all the effects coming from the $`O(k^{\mathrm{\hspace{0.17em}2}j})`$ terms of the $`\stackrel{~}{G}_{2,\mathrm{\hspace{0.17em}2}j}^{(m=\mathrm{\hspace{0.17em}0})}(\stackrel{}{k})`$’s is seen to be in the same order of density of scatterers as the $`O(k^0)`$ term on the r.h.s. of Eq. (54). In fact, it also turns out that the $`O(k^{\mathrm{\hspace{0.17em}2}j})`$ terms of the $`\stackrel{~}{G}_{2,\mathrm{\hspace{0.17em}2}j}^{(m=\mathrm{\hspace{0.17em}0})}(\stackrel{}{k})`$’s are the only ones among the $`\stackrel{~}{G}_{2,j}^{(m)}(\stackrel{}{k})`$’s that contribute to $`f_1^{\text{(R)}}`$ in the order of $`\epsilon ^2`$. This implies that along with the solutions, Eqs. (51), (53) and (54), we also need to include the $`O(k^{\mathrm{\hspace{0.17em}2}j})`$ term of $`\stackrel{~}{G}_{2,\mathrm{\hspace{0.17em}2}j}^{(m=\mathrm{\hspace{0.17em}0})}(\stackrel{}{k})`$ to be consistent. If one just considers this $`O(k^{\mathrm{\hspace{0.17em}2}j})`$ term in $`\stackrel{~}{G}_{2,\mathrm{\hspace{0.17em}2}j}^{(m=\mathrm{\hspace{0.17em}0})}(\stackrel{}{k})`$, then it is easy to see that they satisfy a recurrence relation for $`j1`$ :
$`\stackrel{~}{G}_{2,\mathrm{\hspace{0.17em}2}(j+\mathrm{\hspace{0.17em}1})}^{(m=\mathrm{\hspace{0.17em}0})}(\stackrel{}{k})`$ $`=`$ $`{\displaystyle \frac{k_x^2}{v^2k^4}}\stackrel{~}{G}_{2,\mathrm{\hspace{0.17em}2}j}^{(m=\mathrm{\hspace{0.17em}0})}(\stackrel{}{k}),`$ (55)
i.e,
$`\stackrel{~}{G}_{2,\mathrm{\hspace{0.17em}2}j}^{(m=\mathrm{\hspace{0.17em}0})}(\stackrel{}{k})`$ $`=`$ $`{\displaystyle \frac{(1)^{j\mathrm{\hspace{0.17em}1}}}{\sqrt{2\pi V}}}\left({\displaystyle \frac{k_x^2}{v^2k^4}}\right)^j.`$ (56)
The solutions, Eqs. (51), (53), (54) and (56), are then used to determine the integration constants that arise when we solve the differential Eqs. (44-46). It is important to note that the first term on the right hand side of Eq. (54) is inversely proportional to $`k^2`$. This is the origin of the logarithmic terms we find below.
### 4.2 Solution for $`g_2`$
Here we apply the same procedure to solve for the $`\stackrel{~}{g}_2(\stackrel{}{k})`$’s from the equations $`L_\stackrel{}{k}\stackrel{~}{g}_2(\stackrel{}{k})=b_j(\stackrel{}{k})`$ for $`j=\mathrm{\hspace{0.17em}0},\mathrm{\hspace{0.17em}1},\mathrm{\hspace{0.17em}2}`$. This time, the elements of $`L_\stackrel{}{k}`$ are differential operators in the variable $`\rho `$ and the corresponding constants of integrations are determined using the solutions of $`\stackrel{~}{G}_2^{(m)}(\stackrel{}{k})`$’s while maintaining that $`\stackrel{~}{g}_2^{(m)}(\stackrel{}{k})`$’s go to zero as $`\rho 0`$ and as $`\rho \mathrm{}`$. We also note that for our purpose, solutions of the $`\stackrel{~}{g}_2(\stackrel{}{k})`$’s are only needed for $`\rho >{\displaystyle \frac{a}{2}}`$ as the solution of the $`\stackrel{~}{g}_2(\stackrel{}{k})`$’s for $`\rho <{\displaystyle \frac{a}{2}}`$ gives rise to higher order density corrections than under consideration here. These solutions can also be obtained by the mode expansion technique discussed above in the paragraph preceding Section 4.1. However, as it is fairly straightforward to solve the differential Eqs. (44-46) for $`\rho >{\displaystyle \frac{a}{2}}`$, we directly write down the necessary solutions up to $`O(k^0)`$.
We obtain, for $`\rho >{\displaystyle \frac{a}{2}}`$,
$`\stackrel{~}{g}_{2,\mathrm{\hspace{0.17em}0}}^{(m)}(\stackrel{}{k})={\displaystyle \frac{2na}{\sqrt{2\pi V}}}(1\mathrm{\hspace{0.17em}2}na\rho )e^{\mathrm{\hspace{0.17em}2}na\rho },`$
$`\stackrel{~}{g}_{2,\mathrm{\hspace{0.17em}1}}^{(m=1)}(\stackrel{}{k})={\displaystyle \frac{1}{v\sqrt{2\pi V}}}[{\displaystyle \frac{k_x(k_x+ik_y)}{8k^2}}+{\displaystyle \frac{k_x(k_x+ik_y)}{2k^2}}\mathrm{\hspace{0.17em}2}na\rho +{\displaystyle \frac{2na\rho }{8}}`$
$`+{\displaystyle \frac{(2na\rho )^2}{2}}{\displaystyle \frac{(2na\rho )^3}{4}}]e^{\mathrm{\hspace{0.17em}2}na\rho },`$
$`\stackrel{~}{g}_{2,\mathrm{\hspace{0.17em}1}}^{(m=\mathrm{\hspace{0.17em}0})}(\stackrel{}{k})={\displaystyle \frac{ik_x}{vk^2\sqrt{2\pi V}}}\mathrm{\hspace{0.17em}2}nae^{\mathrm{\hspace{0.17em}2}na\rho },`$
$`\stackrel{~}{g}_{2,\mathrm{\hspace{0.17em}1}}^{(m=\mathrm{\hspace{0.17em}1})}(\stackrel{}{k})={\displaystyle \frac{1}{v\sqrt{2\pi V}}}[{\displaystyle \frac{k_x(k_xik_y)}{8k^2}}+{\displaystyle \frac{k_x(k_xik_y)}{2k^2}}\mathrm{\hspace{0.17em}2}na\rho +{\displaystyle \frac{2na\rho }{8}}`$
$`+{\displaystyle \frac{(2na\rho )^2}{2}}{\displaystyle \frac{(2na\rho )^3}{4}}]e^{\mathrm{\hspace{0.17em}2}na\rho },`$
$`\stackrel{~}{g}_{2,\mathrm{\hspace{0.17em}2}}^{(m=\mathrm{\hspace{0.17em}0})}(\stackrel{}{k})={\displaystyle \frac{k_x^2}{v^2k^4\sqrt{2\pi V}}}\mathrm{\hspace{0.17em}2}nae^{\mathrm{\hspace{0.17em}2}na\rho }`$ (57)
$`+{\displaystyle \frac{1}{2v^2\sqrt{2\pi V}}}[{\displaystyle \frac{45}{128na}}{\displaystyle \frac{315k_y^2}{256nak^2}}+{\displaystyle \frac{k_x^2}{4k^2}}\rho {\displaystyle \frac{11}{8}}\mathrm{\hspace{0.17em}2}na\rho ^2`$
$`{\displaystyle \frac{13}{8}}{\displaystyle \frac{k_x^2}{k^2}}\mathrm{\hspace{0.17em}2}na\rho ^2+{\displaystyle \frac{19}{24}}(2na)^2\rho ^3+{\displaystyle \frac{k_x^2}{2k^2}}(2na)^2\rho ^3`$
$`+{\displaystyle \frac{13}{24}}(2na)^3\rho ^4{\displaystyle \frac{1}{8}}(2na)^4\rho ^5]e^{2na\rho }`$
and for the $`O(k^{\mathrm{\hspace{0.17em}2}j})`$ terms in $`\stackrel{~}{G}_{2,\mathrm{\hspace{0.17em}2}j}^{(m=\mathrm{\hspace{0.17em}0})}(\stackrel{}{k})`$ we have
$`\stackrel{~}{g}_{2,\mathrm{\hspace{0.17em}2}j}^{(m=\mathrm{\hspace{0.17em}0})}(\stackrel{}{k})`$ $`=`$ $`{\displaystyle \frac{(1)^{j\mathrm{\hspace{0.17em}1}}}{\sqrt{2\pi V}}}\left({\displaystyle \frac{k_x^2}{v^2k^4}}\right)^j\mathrm{\hspace{0.17em}2}nae^{\mathrm{\hspace{0.17em}2}na\rho }.`$ (58)
We point out that all of the terms in each of the square brackets, in each of the above three equations, are of the same order in the density. This can be seen easily by noting that $`\rho `$ is typically of order $`(2na)^1`$, so that $`(2na\rho )`$ is typically independent of the density. The solutions, Eqs. (51), (53), (54), (56), (57) and (58) now can be assembled to calculate $`G_2`$ and $`g_2`$ in $`(\stackrel{}{r},\stackrel{}{v})`$ and $`(\stackrel{}{r},\stackrel{}{v},\rho )`$ space respectively and feed the results into the r.h.s. of Eqs. (31) and (36) to obtain $`f_1^{\text{(R)}}`$. This involves a summation of different $`m`$ and $`\stackrel{}{k}`$-values. In the infinite volume limit the $`\stackrel{}{k}`$-sum can be converted to an integration over $`\stackrel{}{k}`$. The sum over $`m`$ is straightforward, but we have to remember that the integration over $`\stackrel{}{k}`$ has to be carried out in a range $`kk_0l^1`$. Secondly, since we have expanded the distribution functions in powers of $`\epsilon `$ and then subsequently in powers of $`k`$, the lower limit of $`k`$ for the $`\stackrel{}{k}`$-integration cannot be taken to be zero. To determine this lower limit of $`k`$ for the $`k`$-integration, we observe that the expansion in $`\epsilon `$ cannot be carried out for those values of $`k`$ where $`k<{\displaystyle \frac{\epsilon }{2v}}`$, so that the value $`{\displaystyle \frac{\epsilon }{2v}}`$ forms a natural lower cut-off for the Fourier transform. Our solutions of $`G_2`$ and $`g_2`$ therefore do not hold for $`k<{\displaystyle \frac{\epsilon }{2v}}`$ and to do a satisfactory perturbation theory in the range $`k<{\displaystyle \frac{\epsilon }{2v}}`$, one needs to consider both the $`\epsilon `$ and $`\stackrel{}{k}`$-dependent terms together. After doing so, one finds that such a perturbation theory does not affect our results at the present density order .
Before performing the integration over $`\stackrel{}{k}`$, we notice that in two dimensions, the numerator of the $`\stackrel{}{k}`$-integral is proportional to $`kdk`$. This means that any part of the solutions of $`\stackrel{~}{g}_2(\stackrel{}{k})`$ or $`\stackrel{~}{G}_2(\stackrel{}{k})`$ having a leading power of $`k`$ of order 2 or higher in the denominator gives rise to a singularity at $`k0`$ for the $`\stackrel{}{k}`$-integral. First, the highest leading power of $`k`$ in the denominators of Eqs. (51), (53), (54) and (57) is $`k^2`$, occurring in $`\stackrel{~}{G}_{2,\mathrm{\hspace{0.17em}2}}^{(m=\mathrm{\hspace{0.17em}0})}(\stackrel{}{k})`$ and $`\stackrel{~}{g}_{2,\mathrm{\hspace{0.17em}2}}^{(m=\mathrm{\hspace{0.17em}0})}(\stackrel{}{k})`$ respectively. These terms proportional to $`k^{\mathrm{\hspace{0.17em}2}}`$ give rise to a logarithmic electric field dependence once the $`\stackrel{}{k}`$-integration is performed for $`{\displaystyle \frac{\epsilon }{2v}}kk_0`$. The rest of the terms in these solutions supply only analytic field dependences that can be expressed as power series in $`\epsilon `$. Secondly, even though the solutions given in Eqs. (56) and (58) have higher powers of $`k`$ than $`k^2`$ in the denominators, they also come with subsequently higher powers of $`\epsilon `$ in their numerators. Thus, when the $`\stackrel{}{k}`$-integration is performed, they contribute terms proportional to $`\epsilon ^2`$ or higher, to $`g_2`$ or $`G_2`$. Consequently, in addition to analytic field dependent terms, in our present approximation we have only one non-analytic field dependent term appearing in $`g_2`$ or $`G_2`$ and that is proportional to $`\stackrel{~}{\epsilon }^2\mathrm{ln}\stackrel{~}{\epsilon }`$. No doubt there exist further non-analytic terms in higher orders in $`\stackrel{~}{\epsilon }`$, but their calculation would require a careful consideration of various terms we have neglected here, such as the repeated ring contributions.
### 4.3 Solution for $`f_1^{\text{(R)}}`$ and the calculation of $`\lambda _+^{\text{(R)}}`$
Once the solutions, Eqs. (51), (53), (54), (56), (57) and (58) are inserted in Eqs. (31) and (37) and the $`\stackrel{}{k}`$-integration is performed in the range $`{\displaystyle \frac{\epsilon }{2v}}<k<k_0`$, we get, by the method described in (40-43), the following equations to be solved to obtain $`f_1^{\text{(R)}}`$ and $`F_1^{\text{(R)}}`$ respectively :
$`\epsilon {\displaystyle \frac{}{\theta }}(\mathrm{sin}\theta f_1^{\text{(R)}})+{\displaystyle \frac{}{\rho }}\left\{\left(v+\rho \epsilon \mathrm{cos}\theta +{\displaystyle \frac{\rho ^2\epsilon ^2\mathrm{sin}^2\theta }{v}}\right)f_1^{\text{(R)}}\right\}+\mathrm{\hspace{0.17em}2}navf_1^{\text{(R)}}`$ (59)
$`=\mathrm{\Theta }({\displaystyle \frac{a}{2}}\rho ){\displaystyle \frac{4v\rho }{a\sqrt{1\left(\frac{2\rho }{a}\right)^2}}}\times `$
$`\times {\displaystyle _{\frac{\pi }{2}}^{\frac{\pi }{2}}}d\varphi \mathrm{cos}\varphi b_\sigma [{\displaystyle \frac{\epsilon ^2}{8\pi ^2v^2}}\left\{\mathrm{ln}{\displaystyle \frac{2vk_0}{\epsilon }}\right\}{\displaystyle \frac{k_0^2}{8\pi ^2}}\{{\displaystyle \frac{3}{16nav}}\epsilon \mathrm{cos}\theta +{\displaystyle \frac{135}{2048(nav)^2}}\epsilon ^2\}+{\displaystyle \frac{A\epsilon ^2}{16\pi ^3v^2}}]`$
$`\mathrm{\hspace{0.17em}2}av[{\displaystyle \frac{\epsilon ^2}{8\pi ^2v^2}}\left\{\mathrm{ln}{\displaystyle \frac{2vk_0}{\epsilon }}\right\}\mathrm{\hspace{0.17em}2}nae^{2na\rho }`$
$`+{\displaystyle \frac{k_0^2}{8\pi ^2}}e^{2na\rho }\{2na(1\mathrm{\hspace{0.17em}2}na\rho ){\displaystyle \frac{\epsilon }{v}}[{\displaystyle \frac{(2na\rho )^3}{2}}(2na\rho )^2{\displaystyle \frac{3}{4}}\mathrm{\hspace{0.17em}2}na\rho +{\displaystyle \frac{1}{8}}]\mathrm{cos}\theta `$
$`{\displaystyle \frac{\epsilon ^2}{4nav^2}}[{\displaystyle \frac{(2na\rho )^5}{8}}{\displaystyle \frac{13}{24}}(2na\rho )^4{\displaystyle \frac{25}{24}}(2na\rho )^3+{\displaystyle \frac{35}{16}}(2na\rho )^2{\displaystyle \frac{2na\rho }{8}}+{\displaystyle \frac{135}{256}}]\}`$
$`+{\displaystyle \frac{1}{16\pi ^3v^2}}A\epsilon ^2\mathrm{\hspace{0.17em}2}nae^{\mathrm{\hspace{0.17em}2}na\rho }]+....`$
and
$`\epsilon {\displaystyle \frac{}{\theta }}(\mathrm{sin}\theta F_1^{\text{(R)}})nav{\displaystyle _{\frac{\pi }{2}}^{\frac{\pi }{2}}}𝑑\varphi \mathrm{cos}\varphi (b_\sigma \mathrm{\hspace{0.17em}1})F_1^{\text{(R)}}`$ (60)
$`={\displaystyle _{\frac{\pi }{2}}^{\frac{\pi }{2}}}d\varphi \mathrm{cos}\varphi (b_\sigma \mathrm{\hspace{0.17em}1})[{\displaystyle \frac{a\epsilon ^2}{8\pi ^2v}}\mathrm{ln}\left({\displaystyle \frac{2vk_0}{\epsilon }}\right)`$
$`{\displaystyle \frac{avk_0^2}{8\pi ^2}}\left\{{\displaystyle \frac{3}{16nav}}\epsilon \mathrm{cos}\theta +{\displaystyle \frac{135}{2048(nav)^2}}\epsilon ^2\right\}`$
$`+{\displaystyle \frac{a}{16\pi ^3v}}A\epsilon ^2+....]`$
$`={\displaystyle \frac{ak_0^2}{16\pi ^2na}}\epsilon \mathrm{cos}\theta +....,`$
where $`b_\sigma `$ has been defined in Eq. (33). The $`A`$-dependent terms in Eqs. (59) and (60) originate from Eqs. (58) and (56) respectively after the $`\stackrel{}{k}`$-integration is carried out. Here $`A`$ is the integral<sup>2</sup><sup>2</sup>2We thank the referee for pointing out an error in a previous calculation of this integral.
$`A={\displaystyle _0^{2\pi }}𝑑\varphi \mathrm{cos}^2\varphi \mathrm{ln}\left[\mathrm{\hspace{0.17em}1}+{\displaystyle \frac{1}{4}}\mathrm{cos}^2\varphi \right]=\mathrm{\hspace{0.17em}0.53536}....`$ (61)
The dominant effect of the ring term on the single particle distribution function, i.e, $`f_1^{\text{(R)}}`$, can now be determined from Eqs. (59) and (60). It is also of some interest to give a crude estimate of the terms that we have neglected. One knows from other studies in the kinetic theory of gases that excluded volume corrections to Boltzmann equation results are the numerically most important corrections, until the density of the system becomes high enough that the mean free path of a particle is less than the size of the particle itself. These excluded volume corrections are provided by the Enskog theory, and this theory can be applied to the Lorentz gas, as well . In our case, the Enskog corrections can be included by replacing the density parameter $`n`$ by $`n(1\pi na^2)^{\mathrm{\hspace{0.17em}1}}n(1+\pi na^2)`$ in the Boltzmann equation. The Enskog correction affects both $`\lambda _0`$ and the $`\epsilon `$-dependent terms in the expressions for $`\lambda _\pm `$ in Eq. (14). Along with the Enskog correction there are other correction terms that affect both $`\lambda _0`$ and the field-dependent terms $`\lambda _\pm `$ at the same density order as the Enskog correction. Also, the terms that have been dropped to obtain Eq. (16) from Eq. (11), contribute to $`\lambda _\pm `$ at the same density order as the Enskog correction. However, since the principal objective of this paper is to investigate the non-analytic contribution of the ring term to the Lyapunov exponents, we will ignore the Enskog and related corrections from our consideration. Thus, using Eqs. (59) and (60), one can express the full solutions of $`f_1`$ and $`F_1`$ as sums of a solution in the Boltzmann regime, a correction due to the ring term and a correction due to the Enskog term, plus all of the other terms we have neglected, as
$`f_1=f_1^{\text{(B)}}+f_1^{\text{(R)}}+....`$ (62)
$`F_1=F_1^{\text{(B)}}+F_1^{\text{(R)}}+....`$ (63)
Consequently, for the positive Lyapunov exponent $`\lambda _+`$ we have,
$`\lambda _+=\lambda _+^{\text{(B)}}+\lambda _+^{\text{(R)}}+....`$ (64)
The solution of $`F_1^{\text{(R)}}`$ is quite straightforward,
$`F_1^{\text{(R)}}={\displaystyle \frac{3ak_0^2}{128\pi ^2(na)^2v}}\epsilon \mathrm{cos}\theta +....`$ (65)
However, to solve for $`f_1^{\text{(R)}}`$ we find that in addition to the analytic field-dependent terms which can be expressed as a power series in $`\epsilon `$, there is a non-analytic field-dependent term in $`f_1^{\text{(R)}}`$ proportional to $`\stackrel{~}{\epsilon }^2\mathrm{ln}\stackrel{~}{\epsilon }`$. Thus, with
$`f_1^{\text{(R)}}=f_{1,\text{analytic}}^{\text{(R)}}+f_{1,\text{non-analytic}}^{\text{(R)}},`$ (66)
we have
$`f_{1,\text{analytic}}^{\text{(R)}}={\displaystyle \frac{ak_0^2}{4\pi ^2}}[\mathrm{\hspace{0.17em}2}na\rho {\displaystyle \frac{(2na\rho )^2}{2}}{\displaystyle \frac{\epsilon \mathrm{cos}\theta }{2nav}}\{{\displaystyle \frac{(2na\rho )^4}{4}}(2na\rho )^3+{\displaystyle \frac{(2na\rho )^2}{8}}+{\displaystyle \frac{2na\rho }{8}}\}`$ (67)
$`+{\displaystyle \frac{\epsilon ^2}{4(nav)^2}}\{{\displaystyle \frac{(2na\rho )^6}{32}}+{\displaystyle \frac{11}{48}}(2na\rho )^5{\displaystyle \frac{(2na\rho )^4}{96}}{\displaystyle \frac{79}{96}}(2na\rho )^3`$
$`+{\displaystyle \frac{3}{32}}(2na\rho )^2{\displaystyle \frac{135}{512}}(2na\rho )+{\displaystyle \frac{135}{512}}\}]e^{2na\rho }`$
$`+{\displaystyle \frac{a\epsilon ^2}{(2\pi )^3v^2}}A(1\mathrm{\hspace{0.17em}2}na\rho )e^{2na\rho }+....\text{for }\rho >{\displaystyle \frac{a}{2}}`$
$`=[{\displaystyle \frac{aA}{(2\pi )^3v^2}}{\displaystyle \frac{135ak_0^2}{512(4\pi nav)^2}}]\{1\sqrt{1\left({\displaystyle \frac{2\rho }{a}}\right)^2}\}\epsilon ^2+....\text{for }\rho <{\displaystyle \frac{a}{2}}`$
and
$`f_{1,\text{non-analytic}}^{\text{(R)}}={\displaystyle \frac{a\epsilon ^2}{4\pi ^2v^2}}\left\{\mathrm{ln}{\displaystyle \frac{2vk_0}{\epsilon }}\right\}(1\mathrm{\hspace{0.17em}2}na\rho )e^{2na\rho }\text{for }\rho >{\displaystyle \frac{a}{2}}`$ (68)
$`={\displaystyle \frac{a\epsilon ^2}{4\pi ^2v^2}}\left\{\mathrm{ln}{\displaystyle \frac{2vk_0}{\epsilon }}\right\}\left[\mathrm{\hspace{0.17em}1}\sqrt{1\left({\displaystyle \frac{2\rho }{a}}\right)^2}\right]\text{for }\rho <{\displaystyle \frac{a}{2}}.`$
Notice that the ring contribution to the distribution function in Eqs. (67) and (68) satisfies the boundary conditions that $`f_1^{\text{(R)}}0`$ as $`\rho 0`$ and $`\rho \mathrm{}`$. Equations (67) and (68) also satisfy continuity at $`\rho ={\displaystyle \frac{a}{2}}`$ at the leading density order. The distribution functions, Eqs. (65), (67) and (68), are all the ones that we need to calculate $`\lambda _+^{\text{(R)}}`$. Consequently,
$`\lambda _+^{\text{(R)}}=\lambda _{+,\text{analytic}}^{\text{(R)}}+\lambda _{+,\text{non-analytic}}^{\text{(R)}}`$ (69)
and using the the definition of Lyapunov exponents in Eq. (12), we have
$`\lambda _{+,\text{analytic}}^{\text{(R)}}={\displaystyle _0^{2\pi }}𝑑\theta {\displaystyle _{\frac{a}{2}}^{\mathrm{}}}𝑑\rho {\displaystyle \frac{f_{1,\text{analytic}}^{\text{(R)}}}{\rho }}`$ (70)
$`={\displaystyle \frac{ak_0^2v}{4\pi }}{\displaystyle \frac{ak_0^2l^2\epsilon ^2}{2\pi v}}\left\{{\displaystyle \frac{13}{96}}{\displaystyle \frac{135}{512}}\left(\mathrm{ln}2na^2+𝒞\right)\right\}`$
$`\mathrm{\hspace{0.17em}0.53536}{\displaystyle \frac{a\epsilon ^2}{(2\pi )^2v}}\left(\mathrm{ln}\mathrm{\hspace{0.17em}2}na^2+𝒞\right)+....`$
and
$`\lambda _{+,\text{non-analytic}}^{\text{(R)}}={\displaystyle _0^{2\pi }}𝑑\theta {\displaystyle _{\frac{a}{2}}^{\mathrm{}}}𝑑\rho {\displaystyle \frac{f_{1,\text{ }\text{non-analytic}}^{\text{(R)}}}{\rho }}`$ (71)
$`={\displaystyle \frac{a\epsilon ^2}{2\pi v}}\left\{\mathrm{ln}{\displaystyle \frac{2k_0v}{\epsilon }}\right\}\left(\mathrm{ln}2na^2+𝒞\right),`$
where $`l`$ is the mean free path and $`𝒞`$ is Euler’s constant, $`𝒞=\mathrm{\hspace{0.17em}0.5772}...`$.
### 4.4 Calculation of $`\lambda _{}^{\text{(R)}}`$
To calculate the corresponding effect of the ring term on $`\lambda _{}`$, we make use of the relation Eq. (13). It is easy to calculate the effect of the ring term on $`<\alpha >_{\text{NESS}}`$ using $`F_1^{\text{(R)}}`$ already determined in the previous section. Thus, using
$`\lambda _++\lambda _{}=<\alpha >_{\text{NESS}},`$ (72)
and a complete analogy to Eqs. (62-64), we can calculate three terms of $`<\alpha >_{\text{NESS}}`$:
$`<\alpha >_{\text{NESS}}=<\alpha >_{\text{NESS}}^{\text{(B)}}+<\alpha >_{\text{NESS}}^{\text{(R)}}+....,`$ (73)
with
$`<\alpha >_{\text{NESS}}^{\text{(R)}}={\displaystyle \frac{3ak_0^2l^2\epsilon ^2}{32\pi v}}+....`$ (74)
Following Eqs. (62-64), we now express $`\lambda _{}`$ as $`\lambda _{}=\lambda _{}^{\text{(B)}}+\lambda _{}^{\text{(R)}}+....`$, satisfying $`\lambda _+^{\text{(R)}}+\lambda _{}^{\text{(R)}}=<\alpha >_{\text{NESS}}^{\text{(R)}}`$. This leads us to
$`\lambda _{,\text{analytic}}^{\text{(R)}}={\displaystyle \frac{ak_0^2v}{4\pi }}{\displaystyle \frac{ak_0^2l^2\epsilon ^2}{2\pi v}}\left\{{\displaystyle \frac{5}{96}}+{\displaystyle \frac{135}{512}}(\mathrm{ln}2na^2+𝒞)\right\}`$ (75)
$`+\mathrm{\hspace{0.17em}\hspace{0.17em}0.53536}{\displaystyle \frac{a\epsilon ^2}{(2\pi )^2v}}(\mathrm{ln}\mathrm{\hspace{0.17em}2}na^2+𝒞)+....,`$
and
$`\lambda _{,\text{non-analytic}}^{\text{(R)}}={\displaystyle \frac{a\epsilon ^2}{2\pi v}}\left\{\mathrm{ln}{\displaystyle \frac{2k_0v}{\epsilon }}\right\}(\mathrm{ln}2na^2+𝒞).`$ (76)
where $`l`$ is the mean free path and $`𝒞`$ is Euler’s constant, $`𝒞=\mathrm{\hspace{0.17em}0.5772}...`$.
## 5 The field-dependent collision frequency and its <br>effects on the Lyapunov exponents
As stated before, our second main purpose was the derivation of the leading non-analyticity in the field dependence of the Lyapunov exponents. In analogy with the transport coefficients, we expected these non-analyticities to result from the long time behavior of the ring terms, which we found confirmed in the preceding section. Some further thought reveals we can estimate the non-analytic field dependence in a simple way.
In the presence of a thermostatted field there are two types of contributions to the positive Lyapunov exponent of the two-dimensional Lorentz gas:
1) contributions from the bending of the trajectories by the fields and
2) contributions from the divergence of trajectory pairs at collisions.
The first type of contributions are of order $`\stackrel{~}{\epsilon }^2`$ in the Boltzmann approximation. We expect that the coefficient of this term will pick up higher density corrections and there will be additional terms of higher orders in $`\stackrel{~}{\epsilon }`$. But we have not found any indications for corrections of lower order than $`\stackrel{~}{\epsilon }^2`$ resulting from the field-bending contributions.
The collisional contributions can be generally expressed as an average of the form $`\nu \mathrm{ln}{\displaystyle \frac{|\delta \stackrel{}{v^{}}|}{|\delta \stackrel{}{v}|}}_\text{c}`$, with $`\delta \stackrel{}{v^{}}`$ and $`\delta \stackrel{}{v}`$ the velocity differences between the adjacent trajectories just after and just before a collision, respectively, $`\nu `$ the average collision frequency, and the angular brackets, $`_\text{c}`$, indicating an average over collisions. At low densities even correlated collisions happen at large distances, i.e. in the order of a mean free path length apart from each other. Therefore their distribution of collision angles and hence their contribution to the average $`_\text{c}`$, to the leading order in density remains the same as for uncorrelated collisions. We should then expect that at low densities the main effect of the correlated collisions on the Lyapunov exponents should be due to a change of the collision frequency $`\nu `$ as a result of correlated collisions taking place in the presence of the field. If the latter changes from $`\nu _0`$ to $`\nu _0+\delta \nu `$, then Eq. (15) predicts a change of the positive Lyapunov exponent of magnitude
$`\delta \lambda _+=\delta \nu \left\{\mathrm{ln}{\displaystyle \frac{a\nu _0}{v}}+𝒞\right\}.`$ (77)
To obtain this result we have used the fact that the equilibrium, low density Lyapunov exponent, Eq. 15) can be written in the form
$`\lambda _0=\nu _0\left\{1𝒞\mathrm{ln}{\displaystyle \frac{a\nu _0}{v}}\right\},`$ (78)
where $`\nu _0=2nav`$. In order to understand why and how the thermostatted field changes the collision frequency we first recall that in equilibrium the collision frequency can be obtained simply by using the uniformity of the equilibrium distribution for the point particle in available phase space, with the result that $`\nu ={\displaystyle \frac{2nav}{1\pi na^2}}`$. One just has to consider the probability that the light particle during an infinitesimal time $`dt`$ will hit one of the scatterers. On the other hand, at a time $`t`$ after a given initial time, the probability for a collision may be considered to be a sum of three contributions: the collision frequency obtained by assuming that all collisions are uncorrelated and independent of each other, plus the probability for a recollision with a scatterer with which it has collided before, minus the reduction of the collision probability due to any collected knowledge of where no scatterers are present. In equilibrium the last two contributions have to cancel, as we demonstrate in the Appendix. In the presence of a field, however, this cancellation does not occur. This can easily be understood in a qualitative way following the argument that the cancellation in equilibrium occurs because the probability for return to the boundary of a scatterer is exactly the same as that for return to the boundary of a region where a scatterer could be, but in fact is not present (a virtual scatterer). In the presence of a field, the average velocity of the point particle before collision with a real scatterer will be in the direction of the field, and after the collision the average velocity will be anti-parallel to the field. The field will then tend to turn the particle around and have it move back in the direction of the scatterer. This effect enhances the probability of a recollision in comparison to that for an isotropic distribution around the scatterer. In a “virtual collision”, in which the velocity does not change, the particle, on average, ends up downstream (i.e. in the direction of the applied field) from the virtual scatterer and its recollision probability is decreased compared to that for an isotropic distribution.
In the Appendix, a quantitative calculation is given based on the following two assumptions:
1) After the real or virtual collision the spatial distribution of the point particle becomes centered around a point at a distance of a diffusion length from the scatterer and
2) for long times this distribution can be found by solving the diffusion equation. The resulting expression for $`\delta \nu `$ is
$`\delta \nu ={\displaystyle \frac{a\epsilon ^2}{2\pi v}}\mathrm{ln}{\displaystyle \frac{\nu _0}{\epsilon }}.`$ (79)
A more formal, but equivalent, way to obtain this result is by extending the method described by Latz, van Beijeren and Dorfman for the low density distribution of time of free flights of the moving particle to include the contribution from ring events, so as to apply to a system in a thermostatted electric field. The main idea is to solve a kinetic equation for $`f(\stackrel{}{r},\stackrel{}{v},t,\tau )`$, the distribution of particles at a phase point $`(\stackrel{}{r},\stackrel{}{v})`$ at time $`t`$ such that their last collision took place at a time $`\tau `$ earlier, i.e., at time $`t\tau `$. It is then easy to argue that the distribution of free flight times is simply the derivative of this (“last collision”) distribution with respect to $`t\tau `$. We can then obtain a NESS average of the time of free flight and thereby calculate the field dependent collision frequency $`\nu (\epsilon )=\nu _0+\delta \nu `$. Since we want to show that the origin of the non-analytic field dependence of both $`\lambda _+`$ and $`\lambda _{}`$ is rooted in the non-analytic field dependence of collision frequency $`\delta \nu `$, let us keep only the non-analytic field-dependent term as the leading term of the expansion of $`\delta \nu `$ in the density of scatterers and in the electric field strength and write
$`\delta \nu =\beta \epsilon ^2\mathrm{ln}\left\{{\displaystyle \frac{2k_0v}{\epsilon }}\right\}+\mathrm{},`$ (80)
where the quantity $`\beta `$ has to be determined from the NESS average of $`\tau `$, using the effect of the ring term on the NESS distribution function $`f(\stackrel{}{r},\stackrel{}{v},\tau )`$ with $`k_0`$ of the order of $`{\displaystyle \frac{1}{\nu _0v}}`$. To obtain this distribution function, we follow exactly the same procedure as outlined in Sections 4 and 5, but this time, with the variable $`\tau `$ instead of $`\rho `$. Notice that, this time, even though the equations for corresponding $`f_1`$ and $`g_2`$’s are different, due to the difference in the dynamical equations for $`\dot{\rho }`$ and $`\dot{\tau }`$ during free flights and at collisions, the equations involving $`F_1`$ and $`G_2`$’s remain the same. The source of the non-analytic field-dependent term will surface again exactly from the $`O(k^2)`$ term in Eq. (54). As far as this non-analytic field-dependent term is concerned, at the lowest order of density, the variables $`\rho `$ and $`\tau `$ are identical up to a multiplicative factor $`v`$. Both grow linearly with time in between collisions and are set back to (for $`\rho `$, almost) zero at each collision with a scatterer. One then recovers the corresponding non-analytic part of the NESS distribution function , analogous to Eq. (68),
$`f_{1,\text{non-analytic}}^{\text{(R)}}(\stackrel{}{v},\tau )={\displaystyle \frac{a\epsilon ^2}{4\pi ^2v}}\left\{\mathrm{ln}{\displaystyle \frac{2vk_0}{\epsilon }}\right\}(1\mathrm{\hspace{0.17em}2}nav\tau )e^{2nav\tau }\text{for }\tau >0,`$ (81)
from which $`\beta `$ can be obtained to be
$`\beta ={\displaystyle \frac{a}{2\pi v}},`$ (82)
after which, one easily recovers the result of Eq. (79).
## 6 Discussion
While much of this paper is quite technical, there are two main points that we would like to emphasize: (1) We have developed a method which allows an extension of the calculation of the Lyapunov exponents for a two-dimensional Lorentz gas to higher densities than is possible by means of the ELBE. (2) The logarithmic terms obtained here, while small, are indicators of similar logarithmic terms which are certain to appear when these calculations are extended to general two-dimensional gases, where all of the particles move.
The first point allows one to contemplate a general kinetic theory for the calculation of of sums, at least, of all positive, or of all negative Lyapunov exponents. Such an approach was also indicated by Dorfman, Latz, and van Beijeren , for the KS-entropy of a dilute gas in equilibrium, but the theory there has not yet been developed beyond the Boltzmann equation. The relevance of the second point can be seen if one realizes that the linear Navier-Stokes transport coefficients of a two-dimensional gas diverge because of long time tail effects, of the type discussed here . In the general gas case therefore the logarithmic terms in the positive and negative Lyapunov exponents will not cancel as they do here, because the transport coefficients themselves should diverge as $`\mathrm{ln}\stackrel{~}{\epsilon }`$ as $`\epsilon `$ approaches zero. Thus the logarithmic terms obtained here should be seen as precursors of the more important logarithmic terms that will appear in the theory of two-dimensional gases.
It is worth noting that the $`\stackrel{~}{\epsilon }^2\mathrm{ln}\stackrel{~}{\epsilon }`$ term results from a long range correlation in time between the moving particle and the scatterers that is present in both the pair correlation functions, $`G_2`$, and $`g_2`$, either of which is proportional to the square of the electric field strength and the inverse square of the wave number, at small wave numbers and fields. This dependence is not present in the Lorentz gas in equilibrium, of course, but similar collision frequency arguments to those given here suggest that non-analytic terms may be present in the ring contributions to the positive Lyapunov exponent for trajectories on the fractal repeller for an open Lorentz gas. In this case the inverse system size, $`L^1`$, plays the role of $`{\displaystyle \frac{\epsilon }{2v}}`$, the lower limit of $`k`$ for the integration over $`\stackrel{}{k}`$ and one would expect to find terms of order $`L^2\mathrm{ln}L`$ in the ring term for this case. This point is currently under investigation.
Finally we mention that neither the non-analytic terms found here, nor the excluded volume corrections included in the Enskog terms are able to account for the field dependence of the Lyapunov exponents as observed in the computer simulations by Dellago and Posch . This is not unexpected since we have not been systematic in computing the density dependence of the coefficient of $`\epsilon ^2`$, nor have we considered higher order terms in $`\epsilon `$ beyond order $`\epsilon ^2\mathrm{ln}\epsilon `$. All of the neglected terms are likely to be numerically more important than the ones we have kept. There is also no indication in the simulation data for the Lyapunov exponents of a clear presence of the interesting logarithmic term in the applied field. Such logarithmic terms are typically difficult to detect in simulation data, without a careful hunt for them. However, it may be easier to check, by means of computer simulation, the existence of the $`\stackrel{~}{\epsilon }^2\mathrm{ln}\stackrel{~}{\epsilon }`$ term in the collision frequency than in the Lyapunov exponents. In any case, we would like to emphasize that computer simulation studies of thermostatted systems provide very useful ways to check a number of phenomena predicted by the kinetic theory of moderately dense gases.
ACKNOWLEDGEMENTS: This paper is dedicated to George Stell on the occasion of his 60-th birthday. We would like to thank Kosei Ide, Zoltán Kovács, Herman Kruis, Arnulf Latz, Luis Nasser, and Ramses van Zon for many valuable conversations and suggestions. J. R. D. wishes to thank the National Science Foundation for support under Grant PHY-96-00428. H. v. B. acknowledges support by FOM, SMC and by the NWO Priority Program Non-Linear Systems, which are financially supported by the ”Nederlandse Organisatie voor Wetenschappelijk Onderzoek (NWO)”.
## Appendix <br>Derivation of the field-dependent collision frequency
To derive the field dependence of the collision frequency we first approximate the probability of a recollision at time $`t`$ as
$`P^{\text{rec}}(t)={\displaystyle \frac{\nu }{2}}{\displaystyle _0^{2\pi }}𝑑\theta {\displaystyle _0^{\mathrm{}}}𝑑\tau {\displaystyle _{\stackrel{}{v}\widehat{\sigma }>0}}𝑑\widehat{\sigma }|\stackrel{}{v}\widehat{\sigma }|R(\tau ,\theta ,\sigma )b_{\widehat{\sigma }}F^{\text{(B)}}(\theta ).`$ (A1)
Here $`F^{\text{(B)}}(\theta )`$ describes the Boltzmann distribution for the velocity in the NESS. The function $`R(\tau ,\theta ,\sigma )`$ describes the probability density for return to the circumference of a given scatterer in a time $`\tau `$ just after colliding with this scatterer with scattering vector $`\widehat{\sigma }`$ and post-collisional velocity described by $`\theta `$ (see Fig. 4). We have ignored a possible dependence of the collision frequency $`\nu `$ on $`\widehat{v}`$, which would only play a role at higher orders in the density.
Fig. 4 : A recollision taking place after a real (solid line) or corresponding virtual (dashed line) collision, followed by a post-collisional excursion maintaining on average the direction of velocity over a persistence length $`l_p`$.
Similarly the reduction of the collision frequency at time $`t`$ due to virtual recollisions can be estimated as
$$P^{\text{nc}}(t)=\frac{\nu }{2}_0^{2\pi }𝑑\theta _0^{\mathrm{}}𝑑\tau _{\stackrel{}{v}\widehat{\sigma }>0}𝑑\widehat{\sigma }|\stackrel{}{v}\widehat{\sigma }|R(\tau ,\theta ,\sigma )F^{\text{(B)}}(\theta ).$$
(A2)
In equilibrium $`F^{\text{(B)}}(\theta )`$ is independent of $`\theta `$, so one sees immediately that both terms cancel, as they should. In the presence of a thermostatted field we need the explicit form of $`F_1^{\text{(B)}}(\theta )`$ up to the first field-dependent order, given in Eq. (50) as
$$F_1^{\text{(B)}}(\theta )=\frac{1}{2\pi }\left[\mathrm{\hspace{0.17em}1}+\frac{3\epsilon }{8nav}\mathrm{cos}\theta \right].$$
(A3)
The function $`R(\tau ,\theta ,\sigma )`$ for large enough $`\tau `$ may be approximated by the product of $`2av`$ (velocity times cross section) and the probability density for finding the point particle at the position of the scatterer. For weak fields the latter may be approximated by the solution of a diffusion equation with a drift velocity $`u\widehat{x}`$ in the $`+x`$-direction and an initial density localized at the position $`l_p\widehat{\theta }`$ with respect to the center of the scatterer. Here $`l_p`$ is the persistence length, that is, the average distance traveled by a point particle in an equilibrium system in the direction of its initial velocity and $`\widehat{\theta }`$ is the unit vector in the direction of the velocity right after the initial collision at $`t\tau `$. The persistence length may be expressed as $`l_p={\displaystyle _0^{\mathrm{}}}𝑑t\widehat{v}\stackrel{}{v}(t)`$. Multiplying this by the constant speed $`v`$ one finds with the aid of the Green-Kubo expression for the diffusion that $`l_p={\displaystyle \frac{2D}{v}}`$ in two dimensions. This assumption for the long time distribution may be understood by imagining that the first few free flights after the initial collision of the particle move it over a distance in the order of a mean free path in the direction of its initial postcollisional velocity before it starts to diffuse by virtue of further collisions with scatterers. Thus for large $`\tau `$ the distribution of the light particle will be centered around the point $`l_p\widehat{\theta }`$ with respect to the center of the scatterer, and the final point, on the surface of the scatterer, may be approximated to be at the center of the scatterer as well, because of low density. These arguments lead to the explicit form for the recollision probability given by
$$R(\tau ,\theta ,\sigma )=\mathrm{\hspace{0.17em}2}av\frac{e^{\frac{[l_p\widehat{\theta }+u\tau \widehat{x}]^2}{4D\tau }}}{4\pi D\tau }.$$
(A4)
Finally we need the explicit form $`u={\displaystyle \frac{3\epsilon v}{8\nu _0}}`$ for the drift velocity to leading order in the density, and the identity that
$$\frac{1}{2}_{\stackrel{}{v}\widehat{\sigma }>0}𝑑\widehat{\sigma }|\stackrel{}{v}\widehat{\sigma }|b_{\widehat{\sigma }}\mathrm{cos}\theta =\frac{v}{3}\mathrm{cos}\theta .$$
(A5)
Then, after expanding
$$e^{\frac{2l_pu\tau \widehat{\theta }\widehat{x}}{4D\tau }}=\mathrm{\hspace{0.17em}1}\frac{l_pu}{2D}\mathrm{cos}\theta +...,$$
(A6)
we can now do all the calculations needed to obtain the leading non-analytic term in the field expansion of the collision frequency. We find that
$`\delta \nu =av\nu {\displaystyle _0^{2\pi }}𝑑\theta {\displaystyle _0^{\mathrm{}}}𝑑\tau {\displaystyle \frac{e^{\frac{[l_p^2+(u\tau )^2]}{4D\tau }}}{4\pi D\tau }}\left[\mathrm{\hspace{0.17em}1}{\displaystyle \frac{l_pu}{2D}}\mathrm{cos}\theta \right]{\displaystyle _{\stackrel{}{v}\widehat{\sigma }>0}}𝑑\widehat{\sigma }|\stackrel{}{v}\widehat{\sigma }|(b_{\widehat{\sigma }}1){\displaystyle \frac{3\epsilon \mathrm{cos}\theta }{16\pi nav}}.`$ (A7)
After performing the integrations, we recover Eq. (79). Notice that the logarithm of $`\stackrel{~}{\epsilon }`$ results from the cut-off on the $`\tau `$ integration provided by the drift term in the exponential. |
warning/0002/cond-mat0002197.html | ar5iv | text | # Anomalous scaling dimensions and stable charged fixed-point of type-II superconductors
## Abstract
The critical properties of a type-II superconductor model are investigated using a dual vortex representation. Computing the propagators of gauge field $`𝐀`$ and dual gauge field $`𝐡`$ in terms of a vortex correlation function, we obtain the values $`\eta _𝐀=1`$ and $`\eta _𝐡=1`$ for their anomalous dimensions. This provides support for a dual description of the Ginzburg-Landau theory of type-II superconductors in the continuum limit, as well as for the existence of a stable charged fixed point of the theory, not in the $`3DXY`$ universality class.
Determining the universality class of the phase-transition in a system of a charged scalar field coupled to a massless gauge field, such as a type-II superconductor, has been a long-standing problem . Analytical and numerical efforts have recently focused on the use of a dual description of the Ginzburg-Landau theory (GLT) of type-II superconductors, pioneered by Kleinert , in investigating the character of a proposed novel stable fixed point of the theory for a charged superconducting condensate, in which case the $`3DXY`$ fixed point of the neutral superfluid is rendered unstable . The dual formulation has also been employed to investigate the possibility of novel broken symmetries in the vortex liquid phase of such systems in magnetic fields .
The GLT is defined by a complex matter field $`\psi `$ coupled to a massless fluctuating gauge field $`𝐀`$ with a Hamiltonian
$`H_{\psi ,𝐀}=m_\psi ^2\left|\psi \right|^2+{\displaystyle \frac{u_\psi }{2}}\left|\psi \right|^4+\left|\left(i2e𝐀\right)\psi \right|^2+`$ (2)
$`{\displaystyle \frac{1}{2}}\left(\times 𝐀\right)^2.`$
Here, $`e`$ is the electron charge, and $`H_{\psi ,𝐀}`$ is invariant under the local gauge-transformation $`\psi \psi \mathrm{exp}(i\theta )`$, $`𝐀𝐀+\theta /2ie`$. The GLT sustains stable topological objects in the form of vortex lines and vortex loops, the latter are the critical fluctuations of the theory . These objects are highly nonlocal in terms of $`\psi `$, but a dual formulation offers a local field theory for them. The continuum dual representation of the topological excitations, (in $`D=3`$ only), consists of a complex matter field $`\varphi `$ coupled to a massive gauge field $`𝐡`$ , with coupling constant given by the dual charge $`e_d`$, and with dual Hamiltonian
$`H_{\varphi ,𝐡}=m_\varphi ^2\left|\varphi \right|^2+{\displaystyle \frac{u_\varphi }{2}}\left|\varphi \right|^4+\left|\left(ie_d𝐡\right)\varphi \right|^2+`$ (4)
$`{\displaystyle \frac{1}{2}}\left(\times 𝐡\right)^2+{\displaystyle \frac{1}{2}}\left(\times 𝐀\right)^2+ie\left(\times 𝐡\right)𝐀.`$
The massiveness of $`𝐡`$ reduces the symmetry to a global $`U(1)`$-invariance. For details on how to obtain this dual Hamiltonian, we refer the reader to the thorough exposition of this presented in the textbook of Kleinert . For $`e0`$ the original GLT in Eq. 2 has a local gauge symmetry, the dual theory in Eq. 4 has a global $`U(1)`$ symmetry. In the limit $`e0`$, $`𝐀`$ decouples from $`\psi `$ in Eq. 2, $`H_\psi `$ describes a neutral superfluid, and the symmetry is reduced to global $`U(1)`$. The dual Hamiltonian $`H_{\varphi ,𝐡}`$ describes a charged superfluid coupled to a massless gauge field $`𝐡`$ with coupling constant $`e_d`$, and the global symmetry is extended to a local gauge symmetry. Hence, when $`e0`$, the dual of a neutral superfluid is isomorphic to a superconductor. Integrating out the $`𝐀`$ field in Eq. 4 produces a mass-term $`e^2𝐡^2/2`$, where an exact renormalization-group equation for the mass of $`𝐡`$ is given by $`e^2/\mathrm{ln}l=e^2`$ . Therefore, when $`e0`$, then $`e^2\mathrm{}`$ as $`l\mathrm{}`$. This supresses the dual gauge field, and the resulting dual theory is a pure $`|\varphi |^4`$-theory. Hence, in the long-wavelength limit, the dual of a superconductor is isomorphic to a neutral superfluid .
In this paper, we obtain the anomalous scaling dimensions $`\eta _𝐀`$ of the gauge field , as well as $`\eta _𝐡`$ of the dual gauge field, not previously considered, directly from large-scale Monte-Carlo simulations. At a $`3DXY`$ critical point,$`\eta _𝐀=\eta _𝐡=0`$. We find that $`(\eta _𝐀=1,\eta _𝐡=0)`$ when $`e0`$, and that $`(\eta _𝐀=0,\eta _𝐡=1)`$, when $`e=0`$. We also contrast the anomalous dimension of the dual mass field $`\varphi `$ at the dual charged (original neutral) and dual neutral (original charged) fixed points, obtaining $`\eta _\varphi =0.24`$ in the former case, and $`\eta _\varphi =0.04`$ in the latter.
A duality transformation, to a set of interacting vortex loops, is performed on the London/Villain approximation to the GLT. In this approximation the partition function is
$`Z(\beta ,e)={\displaystyle }D𝐀D\theta {\displaystyle \underset{\{𝐧\}}{}}\mathrm{exp}[{\displaystyle \underset{𝐱}{}}\{{\displaystyle \frac{1}{2}}(\mathrm{\Delta }\times 𝐀)^2+`$ (6)
$`{\displaystyle \frac{\beta }{2}}(\mathrm{\Delta }\theta e𝐀2\pi 𝐧)^2\}].`$
Here, $`\theta `$ is the local phase of the superconducting order parameter $`\psi `$, while $`𝐧`$ is an integer-valued velocity field (not vortex field) introduced to make the Villain potential $`2\pi `$-periodic. The symbol $`\mathrm{\Delta }`$ denotes a lattice derivative. Amplitude fluctuations are neglected in this approach. The validity of this approximation for $`3D`$ systems, has recently been investigated in detail, both numerically and analytically .
An auxiliary velocity field $`𝐯`$ linearises the kinetic energy. Performing the $`\theta `$-integration constrains $`𝐯`$ to satisfy the condition $`\mathrm{\Delta }𝐯=0`$, explicitly solved by writing $`𝐯=\mathrm{\Delta }\times 𝐡`$, where $`𝐡`$ is forced to integer values by the summation over $`𝐧`$. Introducing an integer-valued vortex field $`𝐦=\mathrm{\Delta }\times 𝐧`$, and using Poisson’s summation formula, we find
$`S(𝐀,𝐡,𝐦)`$ $`=`$ $`{\displaystyle \underset{𝐱}{}}\{2\pi i𝐦𝐡+{\displaystyle \frac{1}{2\beta }}(\mathrm{\Delta }\times 𝐡)^2`$ (8)
$`+ie(\mathrm{\Delta }\times 𝐡)𝐀+{\displaystyle \frac{1}{2}}(\mathrm{\Delta }\times 𝐀)^2\}.`$
Integrating the gauge field in Eq. 8 produces a mass term $`e^2𝐡^2/2`$, giving an effective theory containing the vortex field $`𝐦`$ coupled to a massive gauge field $`𝐡`$
$`Z(\beta ,e)={\displaystyle }D𝐡{\displaystyle \underset{\{𝐦\}}{}}{\displaystyle \underset{𝐱}{}}\delta _{\mathrm{\Delta }𝐦,0}\mathrm{exp}[{\displaystyle \underset{𝐱}{}}\{2\pi i𝐦𝐡+`$ (10)
$`{\displaystyle \frac{e^2}{2}}𝐡^2+{\displaystyle \frac{1}{2\beta }}(\mathrm{\Delta }\times 𝐡)^2\}].`$
The variables $`𝐦`$ in Eq. 10 describe a set of interacting vortices, where the interactions are mediated through the gauge field $`𝐡`$. The variables in Eq. 10 are defined on a lattice which is dual to the lattice from Eq. 6, and the behavior with respect to temperature is inverted in the new variables. The $`\theta `$ field in Eq. 6 describes order, while the $`𝐦`$ field represents the topological excitations of the $`\theta `$ field. These excitations destroy superconducting coherence, and hence quantify disorder .
Integrating out the $`𝐡`$ field in Eq. 10, we obtain the Hamiltonian employed in the present simulations,
$`H(𝐦)`$ $`=`$ $`2\pi ^2J_0{\displaystyle \underset{𝐱_\mathrm{𝟏},𝐱_\mathrm{𝟐}}{}}𝐦(𝐱_\mathrm{𝟏})V(𝐱_\mathrm{𝟏}𝐱_\mathrm{𝟐})𝐦(𝐱_\mathrm{𝟐}),`$ (11)
$`V(𝐱)`$ $`=`$ $`{\displaystyle \underset{𝐪}{}}{\displaystyle \frac{e^{i𝐪𝐱}}{4_\mu \mathrm{sin}^2\left(\frac{q_\mu }{2}\right)+\lambda ^2}}.`$ (12)
In Eq. 12, the charge $`e`$ and lattice-spacing $`a`$ have both been set to unity, and $`\lambda `$ is the bare London penetration depth. At every MC step, we attempt to insert a loop of unit vorticity and random orientation. A new energy is calculated from Eq. 6, and the proposed move is accepted or rejected according to the Metropolis algorithm. This procedure ensures that the vortex lines of the system always form closed loops of random size and shape . In all simulations, a system size of $`40\times 40\times 40`$ was used, and up to $`1.510^5`$ sweeps over the lattice per temperature were used.
To investigate the properties of $`𝐀`$ and $`𝐡`$ at the charged critical point of the original theory, Eq. 2, we have calculated the correlation functions $`𝐀_𝐪𝐀_𝐪`$ and $`𝐡_𝐪𝐡_𝐪`$ in terms of vortex correlations, obtaining
$`𝐀_𝐪𝐀_𝐪`$ $`=`$ $`{\displaystyle \frac{1}{\left|𝐐\right|^2+m_0^2}}\left(1+{\displaystyle \frac{4\pi ^2\beta m_0^2G(𝐪)}{\left|𝐐\right|^2\left(\left|𝐐\right|^2+m_0^2\right)}}\right),`$ (13)
$`𝐡_𝐪𝐡_𝐪`$ $`=`$ $`{\displaystyle \frac{2\beta }{\left|𝐐\right|^2+m_0^2}}\left(1{\displaystyle \frac{2\beta \pi ^2G(𝐪)}{\left|𝐐\right|^2+m_0^2}}\right),`$ (14)
where $`G(𝐪)=𝐦_𝐪𝐦_𝐪`$, $`m_0=\lambda ^1`$ and $`Q_\mu =1e^{i𝐪\widehat{\mu }}`$. All correlation functions have been calculated in the transverse gauge $`𝐀=𝐡=0`$. Both of the fields $`𝐡`$ and $`𝐀`$ are renormalized by vortex fluctuations, albeit in quite different ways.
Invoking the standard form $`\left(q^2+m_{\mathrm{eff}}^2\right)^1`$ for the correlation functions in the immediate vicinity of the critical point in the limit $`q0`$, we find the following expressions for the effective masses,
$`\left(m_{\mathrm{eff}}^𝐀\right)^2`$ $`=`$ $`\underset{q0}{lim}{\displaystyle \frac{m_0^2}{1+4\pi ^2\beta G(𝐪)q^2}},`$ (15)
$`\left(m_{\mathrm{eff}}^𝐡\right)^2`$ $`=`$ $`\underset{q0}{lim}{\displaystyle \frac{m_0^2}{2\beta \left(1\frac{2\pi ^2\beta G(𝐪)}{m_0^2}\right)}}.`$ (16)
When $`e0`$ the correlation function for $`𝐀`$ assumes the form
$$𝐀_𝐪𝐀_𝐪\frac{1}{q^{2\eta _𝐀}}$$
(17)
at the critical point. To determine $`\eta _𝐀`$, we compute the vortex correlator $`G(q)`$. For $`\lambda <<L=40`$ , we expect the following behaviour for $`G(𝐪)`$ in the limit $`q0`$,
$`T<T_c`$ $``$ $`G(𝐪)q^2,`$ (18)
$`T=T_c`$ $``$ $`G(𝐪)q^\eta ,`$ (19)
$`T>T_c`$ $``$ $`G(𝐪)C(T).`$ (20)
When these limiting forms are inserted in Eq. 15, we see that for $`TT_c`$, $`m_{\mathrm{eff}}^𝐀`$ will be finite through the Higgs Mechanism (Meissner effect). For $`TT_c`$ we will have $`m_{\mathrm{eff}}^𝐀=0`$ as in the normal case of a massless photon. Assuming $`G(q)q^\eta `$ precisely at the critical point, it is seen that $`\eta `$ corresponds to $`\eta _𝐀`$ from Eq. 17. We thus identify the scaling power of $`G(𝐪)`$ at the critical point with the anomalous dimension of the massless gauge field $`𝐀`$.
All three limiting forms Eqs. 18-20 are shown in Fig. 1. The gauge field masses $`m_{\mathrm{eff}}^𝐡`$ and $`m_{\mathrm{eff}}^𝐀`$ in Eqs. 15 and 16, are shown in Fig. 2. At the critical point $`G(q)q`$, so that $`\eta _𝐀=1`$. Note that, while $`m_{\mathrm{eff}}^𝐀`$ vanishes at $`T=T_c`$, $`m_{\mathrm{eff}}^𝐡`$ is finite but non-analytic. As a result of the vortex loop blowout, the screening properties of the vortices are dramatically increased, and $`m_{\mathrm{eff}}^𝐡`$ increases sharply.
To find $`\eta _𝐡`$ independently, we consider first the uncharged case $`\lambda \mathrm{}`$, $`m_00`$. First, at an intermediate step in the transformation Eqs. 6 \- 10, the action reads
$$S(\beta ,e)=\underset{𝐱}{}\left\{\frac{1}{2\beta }𝐥^2+ei𝐀𝐥+\frac{1}{2}\left(\times 𝐀\right)^2\right\}.$$
(21)
Here, $`𝐥`$ is an integer field of closed current loops. Setting $`e=0`$ in Eq. 10, the action of the dual Villain model is obtained,
$$\stackrel{~}{S}_V(\beta ,\mathrm{\Gamma })=\underset{𝐱}{}\{2\pi i𝐦𝐡+\frac{1}{2\beta }\left(\mathrm{\Delta }\times 𝐡\right)^2+\frac{\mathrm{\Gamma }}{2}𝐦^2\}.$$
(22)
Here, a term $`\mathrm{\Gamma }𝐦^2/2`$ has been added, and $`\stackrel{~}{S}_V(\beta ,\mathrm{\Gamma })`$ corresponds to the Villain-action in the limit $`\mathrm{\Gamma }0`$. However, it is physically reasonable to propose that the limit $`\mathrm{\Gamma }0`$ is non-singular, since the added term is short-ranged. It should therefore be an irrelevant perturbation, in renormalization group sense, to the long-ranged Biot-Savart interaction governing the fixed point, which is mediated by $`𝐡`$. Rescaling $`𝐡𝐡e/2\pi `$ in Eq. 22, we have $`Z(\beta ,e)=\stackrel{~}{Z}_V(e^2/4\pi ^2,1/2\beta ))`$, leaving Eqs. 21 and 22 interchangeable; $`\eta _𝐡`$ from Eq. 22 should have the same value as $`\eta _𝐀`$ from Eq. 21. The above is demonstrated by our simulations based on Eqs. 11-14, which are independent of the proposed form Eq. 22.
To determine $`\eta _𝐡`$ we study the correlation function $`𝐡_𝐪𝐡_𝐪`$ (Eq. 14) in the limit $`m_00`$. At the uncharged fixed point of the original theory, which is the charged fixed point of the dual theory, we have $`lim_{q0}2\pi \beta ^2G(𝐪)=(1C_2(T))q^2+\mathrm{},q^2C_3(T)q^{2+\eta _𝐡}+\mathrm{}`$, and $`q^2C_4(T)q^4+\mathrm{}`$, for $`T<T_c,T=T_c`$, and $`TT_c`$, respectively. Here, $`C_2(T)`$ corresponds to the helicity modulus (superfluid density) , $`C_3(T)`$ is a critical amplitude, and $`C_4(T)`$ is the inverse of the mass of the dual gauge field for $`TT_c`$. Correspondingly, we have $`lim_{q0}<𝐡_𝐪𝐡_𝐪>=2\beta C_2/q^2,2\beta C_3/q^{2\eta _𝐡}`$, and $`2\beta C_4`$, for $`T<T_c,T=T_c`$, and $`TT_c`$, respectively. Note that $`𝐡`$ is massless for $`T<T_c`$, while it is massive for $`T>T_c`$, the dual system exhibits a “dual Meissner-effect” for $`TT_c`$. At $`T=T_c`$, we have $`q^2𝐡_𝐪𝐡_𝐪C_3(T)q^{\eta 𝐡}`$. A plot of $`q^2𝐡_𝐪𝐡_𝐪`$ is shown in Fig. 3. A linear behaviour at $`T=T_c`$ is found, implying that $`\eta _𝐡=1`$ when $`e=0`$. Since $`\eta _𝐡=1`$ in the uncharged case, this provides further support for the Hamiltonian Eq. 4.
We now set $`e0`$. The gauge field $`𝐡`$ becomes massive via the term $`e^2𝐡^2/2`$, which appears after integrating out the $`𝐀`$ field in Eq. 4. In this case, $`lim_{q0}𝐡_𝐪𝐡_𝐪=2\beta /m_0^2`$ from Eq. 14, and $`𝐡(r)`$ would naively have the trivial scaling dimension $`\left(2d\right)/2`$. However, the mass term offers us a freedom in assigning dimensions to $`e`$ and $`𝐡`$, by introducing renormalization $`Z`$-factors, here $`e^{}=Z_𝐡^{1/2}e`$ and $`𝐡^{}=Z_𝐡^{1/2}𝐡`$.
Prior to integrating out $`𝐀`$ in Eq. 4, the mass appears in the term $`ie\left(\times 𝐡\right)𝐀`$. Integration of the $`\varphi `$ field, partial or complete, can only produce $`(/ie_d𝐡)`$-terms. In particular, this must hold during integration of fast Fourier-modes of the $`\varphi `$ field. Thus, the term $`i(\times 𝐡)𝐀`$ is renormalisation group invariant, i.e. its prefactor must be dimensionless. In terms of scaled fields, at the charged fixed point of the original theory, we have $`𝐀^{}=Z_𝐀^{1/2}𝐀`$, with $`Z_𝐀l^{\eta _𝐀}`$, $`\eta _𝐀=1`$. For $`𝐡`$, we use $`Z_𝐡l^\mathrm{\Delta }`$, where $`\mathrm{\Delta }`$ is not an anomalous scaling dimension ($`𝐡`$ is massive, cf. Fig. 2), but rather a contribution to the engineering dimension of $`𝐡`$. Inserting this into the crossterm $`ie\left(\times 𝐡\right)𝐀`$, we find the scaling dimension $`\left(\eta _𝐀+\mathrm{\Delta }\right)/21`$, which must vanish. This gives the constraint $`\mathrm{\Delta }=1`$ to avoid conflicting results for $`\eta _𝐀`$.
Remarkably, therefore, the scaling dimension of $`𝐡`$ at $`T=T_c`$ is the same in both cases $`m_0=0`$ and $`m_00`$. The results for $`\eta _𝐀`$ and $`\eta _𝐡`$ in the previous paragraphs, are summed up in Table I.
We next consider the distribution of vortex loop sizes in the model Eq. 12, connecting the vortex loop distribution to the anomalous dimension of $`\varphi `$ at $`T_c`$ both for the case $`e=0`$ and $`e0`$. During the simulations, we sample the distribution of loop-sizes $`D(p)`$, where $`p`$ is the perimeter of a loop. This distribution function can be fitted to the form
$$D(p)p^\alpha e^{\beta p\epsilon (T)},$$
(23)
where $`\epsilon (T)`$ is an effective line-tension for the loops. Figures showing the qualitative features of $`D(p)`$ can be found in Ref. . The critical point is characterised by a vanishing line-tension, and close to the critical point we find that $`\epsilon (T)`$ vanishes as $`\epsilon (T)\left|TT_c\right|^{\gamma _\varphi }`$.
The vortex loops are the topological excitations of the GL and 3DXY models, at the same time they are the real-space representation of the Feynman diagrams of the dual field theory. By sampling $`D(p)`$, we obtain information about the dual field $`\varphi `$, particularly $`\gamma _\varphi `$ can be identified as a susceptibility exponent for the $`\varphi `$ field. Using the scaling relation $`\gamma _\varphi =\nu _\varphi \left(2\eta _\varphi \right)`$, and the important observation that even at the charged dual fixed point $`\nu _\varphi =\nu _{3DXY}`$ , this also gives us a value for the anomalous scaling dimension $`\eta _\varphi `$ when we use the value $`\nu _{3DXY}=0.673`$ .
In Ref. the vortex loops of the 3DXY model have been studied meticulously, yielding the value $`\eta _\varphi (0)=0.18\pm 0.07`$. Since the dual of this model is isomorphic to a superconductor, $`\eta _\varphi (0)`$ should be similar to $`\eta _\psi (e)`$ of the original GLT.
We have studied the vortex loop distribution in both the neutral and the charged case. In the former case we find $`\eta _\varphi 0.24`$, in good agreement with Ref. . In the latter case the dual theory has a $`U(1)`$ symmetry, and we would expect to find $`\eta _\varphi =\eta _{3DXY}`$. The exponent $`\eta _{3DXY}`$ has recently been determined with great accuracy to $`\eta _{3DXY}=0.038`$, whereas we find $`\eta _\varphi 0.04`$ which compares well with this value. Fig. 4 shows $`\epsilon (T)`$ for both the charged and uncharged models. It is evident that they belong to two different universality classes.
In the case $`e0`$, which corresponds to the dual neutral case, the inverse $`\varphi `$-propagator is given by $`G^1=q^2+\mathrm{\Sigma }(q)`$, where $`\mathrm{\Sigma }`$ is a self-energy, and $`\mathrm{\Sigma }(q)q^{2\eta }`$ by definition. This gives a leading order behavior $`G1/q^{2\eta }`$ provided $`\eta >0`$, and we find $`\eta =0.04`$ for this case. On the other hand, for the case $`e=0`$, which corresponds to the dual charged case, dual gauge field fluctuations alter the physics, softening the long-wavelength $`\varphi `$ field fluctuations. We obtain $`G^1=q^4+\mathrm{\Sigma }(q)`$, again with $`\mathrm{\Sigma }(q)q^{2\eta }`$, which now gives a leading order behavior $`G1/q^{2\eta }`$, provided $`\eta >2`$. Our result $`\eta =0.24`$ for the case $`e=0`$ (dual charged) is consistent with this, and also with the absolute bounds $`\eta >2D=1`$, in $`D=3`$.
A consequence of the above is that in $`D=3`$ dimensions, $`\lambda \xi ^{(D2)/(2\eta _𝐀)}=\xi `$ at the charged critical point, in contrast to $`\lambda \sqrt{\xi }`$ at the $`3DXY`$ neutral critical point. Since our results have been obtained directly by MC simulations, they are valid beyond all orders in perturbation theory.
This work has been supported by a grant of computing time from Tungregneprogrammet, Norges Forskningsråd. We thank Para//ab for valuable assistance in optimizing our computer codes for use on the Cray Origin 2000, and Zlatko Tešanović for many useful discussions. |
warning/0002/math-ph0002012.html | ar5iv | text | # The inverse spectral problem for surfaces of revolution
## 0 Introduction
This article is concerned with the inverse spectral problem for metrics of revolution on $`S^2`$. We will assume that our metrics are real analytic and belong to a class $`^{}`$ of rotationally invariant metrics which are of ‘simple type’ and which satisfy some generic non-degeneracy conditions (see Definition (0.1)). In particular, we will assume they satisfy the generalized ‘simple length spectrum’ condition that the length functional on the loop space is a clean Bott-Morse function which takes on distinct values on distinct components of its critical set (up to orientation). Denoting by $`Spec(S^2,g)`$ the spectrum of the Laplacian $`\mathrm{\Delta }_g`$, our main result is the following:
Theorem I Spec: $`^{}\mathrm{I}\mathrm{R}^{+\mathrm{I}\mathrm{N}}`$ is 1-1.
Thus, if $`(S^2,g),(S^2,h)`$ are isospectral surfaces of revolution in $`^{}`$, then $`g`$ is isometric to $`h`$. It would be very desirable to strengthen this result by removing the assumption that $`h^{}`$, thereby showing that metrics in $`^{}`$ are spectrally determined within the entire class of analytic metrics on $`S^2`$ with simple length spectra. The only metric on $`S^2`$ presently known to be spectrally determined in this sense is the standard one (which is known to be spectrally determined among all $`C^{\mathrm{}}`$ metrics). A metric $`h`$ satisfying $`Spec(S^2,h)=Spec(S^2,g)`$ for some $`g^{}`$ must have many properties in common with a surface of revolution of simple type; it would be interesting to explore whether it must necessarily be one.
Let us now be more precise about the hypotheses. First, we will assume that there is an effective action of $`S^1`$ by isometries of $`(S^2,g)`$. The two fixed points will be denoted $`N,S`$ and $`(r,\theta )`$ will denote geodesic polar coordinates centered at $`N`$, with $`\theta =0`$ some fixed meridian $`\gamma _M`$ from $`N`$ to $`S`$. The metric may then be written in the form $`g=dr^2+a(r)^2d\theta ^2`$ where $`a:[0,L]\mathrm{I}\mathrm{R}^+`$ is defined by $`a(r)=\frac{1}{2\pi }|S_r(N)|`$, with $`|S_r(N)|`$ the length of the distance circle of radius $`r`$ centered at $`N`$. For any smooth surface of revolution, the function $`a`$ satisfies $`a^{(2p)}(0)=a^{(2p)}(L)=0,a^{}(0)=1,a^{}(L)=1`$ and two such surfaces $`(S^2,g_i)`$ $`(i=1,2)`$ are isometric if and only if $`L_1=L_2`$ and $`a_1(r)=a_2(r)`$ or $`a_1(r)=a_2(Lr).`$ We will then assume that the metrics belong to the following class $``$ of simple analytic surfaces of revolution \[Bl\]:
(0.1) Definition $``$ is the moduli space of metrics of revolution $`(S^2,g)`$ with the properties:
(i) $`g`$ (equivalently $`a`$) is real analytic;
(ii) $`a`$ has precisely one non-degenerate critical point $`r_o(0,L)`$, with $`a^{\prime \prime }(r_o)<0`$, corresponding to an ‘equatorial geodesic’ $`\gamma _E`$;
(iii) the (non-linear) Poincare map $`𝒫_{\gamma _E}`$ for $`\gamma _E`$ is of twist type (cf. §1).
We denote by $`^{}`$ the subset of metrics with ‘simple length spectra’ in the sense above.
Regarding ‘simple length spectra,’ we recall that the closed geodesics of a surface of revolution come in one-parameter families of a common length, filling out invariant torii $`𝒯`$ for the geodesic flow. The canonical involution $`\sigma (x,\xi )=(x,\xi )`$ takes $`𝒯`$ to its ‘time reversal’ $`𝒯`$, and takes the closed geodesics of $`𝒯`$ to their reversals on $`𝒯`$. A closed geodesic and its reversal have the same length, so the length spectrum is automatically double except for the length $`2L`$ of the torus of meridians, which is $`\sigma `$-invariant. The simple length spectrum hypothesis is that up to time reversal, the common lengths of the closed geodesics on distinct torii are distinct (cf. Definition 1.2.2). In fact, it would be sufficient for the proof that the length $`2L`$ of the ‘meridian’ closed geodesics is not the length of closed geodesics on any other torus. In any case, it is not hard to show that $`^{}`$ is residual in $``$ (cf. Proposition 1.2.4).
The condition (iii) actually appears in the proof in the following way: the quadratic coefficient ‘$`\alpha :=h^{\prime \prime }(0)`$’ of the classical Birkhoff normal form of the metric $`|\xi |_g`$ at the torus of meridian geodesics must be non-vanishing (see Definition (1.4.5) for the precise meaning of $`h(\xi )`$). This condition is used in Proposition (4.1.2) and Corollary (4.1.3) to evaluate the wave invariants for the meridian torus.
As will be explained further below, the main purpose of the non-degeneracy and simplicity conditions is to insure that there are global action-angle variables for the geodesic flow. These conditions rule out several types of surfaces of revolution: First, they rule out Zoll surfaces of revolution, which are degenerate in every possible sense. It is indeed unknown at this time whether real analytic Zoll surfaces of revolution are determined by their spectra. They also rule out ‘peanuts of revolution’ (which have hyperbolic waists) and other natural rotational surfaces such as Liouville torii.
To our knowledge, the strongest prior result on the inverse spectral problem for surfaces of revolution is that of Bruning-Heintze \[B.H\]: smooth surfaces of revolution with a mirror symmetry thru the $`xy`$ plane are spectrally determined among metrics of this kind. There are also a number of proofs that a surface of revolution is determined by the joint spectrum of $`\mathrm{\Delta }`$ and of $`/\theta `$, the generator of the rotational symmetry \[Kac\]\[B\]\[Gur\].
The method-of-proof of Bruning-Heintze was based on the observation that the invariant spectrum can be heard from the entire spectrum. Hence by separating variables the problem can be reduced to the inverse spectral problem for 1D even singular Sturm-Liouville operators, which was solved by Gelfand-Levitan and Marchenko.
Our proof of the Main Theorem is based on different kind of method, and the inverse result presented here is hopefully just one illustration of it. We begin with the facts that simple surfaces of revolution are completely integrable on both the classical and quantum levels and that the Laplacian has a global quantum normal form in terms of action operators. We then study the trace of the wave group and prove that from its singularity expansion we can reconstruct the global quantum normal form. Finally we show that this normal form determines the metric.
This approach is suggested by the recent inverse result of Guillemin \[G.1\], which shows that the microlocal normal form of $`\mathrm{\Delta }`$ around each non-degenerate elliptic closed geodesic can be determined from the wave invariants (see also \[Z.1,2\]). However, there is but one non-degenerate closed geodesic on a simple surface of revolution, so the direct application of this inverse result does not take full advantage of the situation. Rather, it is natural to start from the fact that the wave group is completely integrable in the following strong sense: namely, it commutes with an effective action of the torus $`S^1\times S^1`$ by Fourier Integral operators on $`L^2(S^2)`$. That is, there exist global action operators $`\widehat{I}_1,\widehat{I}_2`$ and a polyhomogeneous symbol $`\widehat{H}`$ of degree 1 on $`\mathrm{I}\mathrm{R}^20`$ such that $`\sqrt{\mathrm{\Delta }}=\widehat{H}(\widehat{I_1},\widehat{I_2})`$. This is the global quantum Birkhoff normal form alluded to above. Our principal tool is the following inverse result:
Main Lemma The wave trace invariants of $`(S^2,g)`$ with $`g^{}`$ determine the quantum normal form $`\widehat{H}.`$
This Lemma does not actually require that $`(S^2,g)`$ be a surface of revolution, but only that the geodesic flow is toric integrable, i.e. commutes with an effective Hamiltonian torus action. It immediately implies that the principal symbol $`H(I_1,I_2)`$ of $`\widehat{H}(\widehat{I}_1,\widehat{I}_2)`$ is a spectral invariant. Since $`H(I_1,I_2)`$ is essentially a global Birkhoff normal form for the metric, the wave invariants determine the symplectic equivalence of the geodesic flow. Thus we have:
Corollary 1 From the wave trace invariants of $`(S^2,g)`$ with $`g^{}`$ we can determine the symplectic equivalence class of $`G^t`$.
Corollary 1 does not however determine the isometry class of a general $`g^{}`$: As will be discussed in §1 (see also \[C.K\]), simple surfaces of revolution are not symplectically rigid unless they are mirror symmetric. Otherwise put, recovery of the classical Birkhoff normal form only determines the even part of the metric in the following sense:
Corollary 2 A metric $`g`$ may be written in the form $`g=[f(cosu)]^2du^2+[sinu]^2d\theta ^2`$ (§1, \[Besse\]). From the wave invariants one can determine the even part of $`f`$. In particular if $`g`$ is mirror symmetric, one can determine $`g`$ among mirror symmetric metrics in $`^{}`$ from its spectrum.
It is interesting to note that symplectic rigidity in the mirror-symmetric case gives a new and self-contained proof of the Bruning-Heintze result, without the use of Marchenko’s inverse spectral theorem for singular Sturm Liouville operators. Although this result is superceded by the Theorem, it may have some future relevance to other inverse problems.
To complete the proof of the Theorem, we therefore have to study the subprincipal terms in $`\widehat{H}`$. The result is:
Final Lemma From $`\widehat{H}`$ one can determine the isometry class of $`g`$.
It is in this last Lemma that we use in full that $`(S^2,g)`$ is a surface of revolution rather than just a surface with toric integrable geodesic flow. We also use in full that $`\widehat{H}`$ is a global quantum normal form rather than a microlocal one at $`\gamma _E`$. In subsequent work we will investigate the analogues of the Final Lemma for the microlocal normal form at $`\gamma _E`$, for more general toric integrable metrics and for metrics isospectral to toric integrable ones.
To close this introduction, we discuss some background and some open problems connected with this work.
The principal motivation for studying the inverse spectral problem for surfaces of revolution is its simplicity. There are to date very few inverse results which determine a metric from its Laplace spectrum within the entire class of metrics, or even within concrete infinite dimensional families. To our knowledge, only the standard $`S^n`$ for $`n6`$ and flat 2-torii are known to be spectrally determined. Hence it is desirable to have a simple model of how an inverse result might go.
A second motivation is a somewhat loose analogy between surfaces of revolution and planar domains. Namely, in both cases the unknown is a function of one variable (the profile curve, resp. the boundary) which completely determines the first return times and angles of geodesics emanating from a transversal. That is, paths of billiard trajectories (broken geodesics) on a bounded planar domain are determined from collisions with the boundary, while paths of geodesics on surfaces of revolution are determined from collisions with a meridian (or with the equator). In analytic cases, it is plausible that the unknown function may be determined in large part by the spectrum of first return times from the local transversal. In the case of an analytic surface of revolution, this length spectrum determines the corresonding Birkhoff normal form and hence by Corollary 2 it determines the even part of the profile curve. Similarly, in the case of an analytic plane domain, it is proved in \[CV.2\] that the even part of the boundary of an analytic domain may be determined from the Birkhoff normal form of the billiard map at a bouncing ball orbit. Hence, there is a similarity in the relation between the unknown function and the local classical Birkhoff normal forms. It is interesting to observe in this context that the rigidity result of Colin de Verdiere \[CV.2\] is quite analogous to the inverse result of Bruning-Heintze.
Some immediate open problems: First, there is the symplectic conjugacy problem mentioned above. Second, can one relax analyticity to smoothness in the Theorem above? This is likely to follow from a more intensive analysis of the wave invariants. Third, can one extend it to other ‘non-simple’ types of surfaces of revolution? The main obstacle is that one will generally not have global action-angle variables or global quantum normal forms. What about completely integrable systems in higher dimensions?
Finally, we would like to thank D.Kosygin for several helpful conversations on this paper, and B.Kleiner for giving us up to date information about the status of the conjugacy problem for geodesic flows on surfaces of revolution.
## 1 Classical dynamics
### 1.1 Global action-angle variables
From a geometric (or dynamical) point of view, the principal virtue of metrics in $``$ is described by the following:
(1.1.1) Proposition Suppose $`g`$ is a real analytic metric of revolution on $`S^2`$ such that $`a`$ has precisely one non-degenerate critical point at some $`r_o(0,L).`$ Then the Hamiltonian $`|\xi |_g:=\sqrt{g^{ij}\xi _i\xi _j}`$ on $`T^{}S^2`$ is completely integrable and possesses global real analytic action-angle variables.
Proof: The complete integrability of $`|\xi |_g`$ (i.e. of the geodesic flow) is classical, and follows from the existence of the Clairaut integral $`p_\theta (v):=v,\frac{}{\theta }`$. Since the Poisson bracket $`\{p_\theta ,|\xi |_g\}=0`$, the geodesics are constrained to lie on the level sets of $`p_\theta `$; and since both $`|\xi |_g`$ and $`p_\theta `$ are homogeneous of degree one, the behaviour of the geodesic flow is determined by its restriction to $`S_g^{}S^2=\{|\xi |_g=1\}`$. With the assumption on $`a`$, the level sets are compact and the only critical level is that of the equatorial geodesics $`\gamma _E^\pm S_g^{}S^2`$ (traversed with either orientation). The other level sets are well-known to consist of two-dimensional torii. (Had we allowed the existence of at least two critical points in $`a`$, there would exist a saddle level, i.e. an embedded non-compact cylinder).
The existence of global action-angle variables follows from the general results of \[D\]\[G.S\] and have been constructed explicitly for simple surfaces of revolution in \[CV.1\]. The general formula is as follows: Let
$$P=(|\xi |_g,p_\theta ):T^{}S^2B:=\{(b_1,b_2):|b_2|a(r_o)b_1\}\mathrm{I}\mathrm{R}\times \mathrm{I}\mathrm{R}^+$$
$`(\mathrm{1.1.2})`$
be the moment map of the Hamiltonian $`\mathrm{I}\mathrm{R}^2`$-action defined by the geodesic flow and by rotation. The singular set of $`P`$ is the closed conic set $`Z:=\{(r_o,\theta ,0,p_\theta ):\theta [0,2\pi ),p_\theta \mathrm{I}\mathrm{R}\}`$, i.e. $`Z`$ is the cone thru the equatorial geodesic (in either orientation). The image of $`Z`$ is the boundary of $`B`$; the map $`P|_{T^{}S_gS^2Z}`$ is a trivial $`S^1\times S^1`$ bundle over the open convex cone $`B_o`$ (the interior of $`B`$). For each $`bB_o`$ , let $`H_1(F_b,ZZ)`$ denote the homology of the fiber $`F_b:=P^1(b).`$ This lattice bundle is trivial since $`B`$ is contractible, so there exists a smoothly varying homology basis $`\{\gamma _1(b),\gamma _2(b)\}H_1(F_b,ZZ)`$ which equals the unit cocircle $`S_N^{}S^2`$ together with the fixed closed meridian $`\gamma _M`$ when $`b`$ is on the center line $`\mathrm{I}\mathrm{R}^+(1,0)`$. The action variables are given by \[CV.1, §6\]
$$I_1(b)=_{\gamma _1(b)}\xi 𝑑x=p_\theta ,I_2(b)=_{\gamma _2(b)}\xi 𝑑x=\frac{1}{\pi }_{r_{}(b)}^{r_+(b)}\sqrt{b_1^2\frac{b_2^2}{a(r)^2}}𝑑r+|b_2|$$
$`(\mathrm{1.1.3})`$
where $`r_\pm (b)`$ are the extremal values of $`r`$ on the annulus $`\pi (F_b)`$ (with $`\pi :S_g^{}S^2S^2`$ the standard projection). On the torus of meridians in $`S_g^{}S^2`$, the value of $`I_2`$ equals $`\frac{L}{\pi }`$ and it equals one on the equatorial geodesic. So extended, $`I_1,I_2`$ are smooth homogeneous functions of degree 1 on $`T^{}S^2`$, and generate $`2\pi `$-periodic Hamilton flows.
It follows that the pair $`:=(I_1,I_2)`$ generates a global Hamiltonian torus ($`S^1\times S^1`$)-action commuting with the geodesic flow. The singular set of $``$ equals $`𝒵:=\{I_2=\pm p_\theta \}`$, corresponding to the equatorial geodesics. The map
$$:T^{}S^2𝒵\mathrm{\Gamma }_o:=\{(x,y)\mathrm{I}\mathrm{R}\times \mathrm{I}\mathrm{R}^+:|x|<y\}$$
$`(\mathrm{1.1.4})`$
is a trivial torus fibration. Henceforth we often write $`T_I`$ for the torus $`^1(I)`$ with $`I\mathrm{\Gamma }_o`$ and let $`\mathrm{\Gamma }`$ denote the closure of $`\mathrm{\Gamma }_o`$ as a convex cone. The symplectically dual angle variables $`(\varphi _1=\theta ,\varphi _2)`$ then give, by definition, the flow times (mod $`2\pi `$) along the orbits of $`(I_1,I_2)`$ from a fixed point on $`F_b`$, which we may take to be the unique point lying above the intersection of the equator and the fixed meridian on $`F_b`$ with the geodesic pointing into the northern hemisphere.
So far, we have only assumed the metric to be $`C^{\mathrm{}}`$. We now observe that if $`g`$ is real analytic, then so are $`I_1,I_2.`$ This is obvious in the case of $`I_1`$ and follows from the explicit formula (1.1.3) for $`I_2.`$
Since the metric norm function $`|\xi |_g`$ commutes with $`I_1,I_2`$, it may be expressed as a function $`H(I)`$ of the action variables. Hamilton’s equations for the geodesic flow then take the form
$$\dot{I}_k=0,\dot{\varphi }_k=\omega _k(I),(k=1,2)$$
where
$$\omega (I)=_IH(I)$$
$`(\mathrm{1.1.5})`$
is the frequency vector of the torus $`T_I`$ with action coordinates $`I`$. The geodesic flow on $`T_I`$ is then given by
$$G^t(I,\varphi )=(I,\varphi +t\omega _I)$$
$`(\mathrm{1.1.5}a)`$
so that all the geodesics are quasi-periodic in action-angle coordinates.
The frequency vector $`\omega _I`$ is homogeneous of degree 0 on $`T^{}S^20`$, and hence is constant on rays of torii $`\mathrm{I}\mathrm{R}^+T_I`$. To break the $`\mathrm{I}\mathrm{R}^+`$ symmetry we restrict to the level set $`\{H(I)=1\}\mathrm{\Gamma }_o`$ in action space and view the frequency vector as the map:
$$\omega :\{H=1\}\mathrm{I}\mathrm{R}^2.$$
Since $`_IH(I)T_I(\{H=1\})`$ the frequency map is more or less the Gauss map of $`\{H=1\}`$ (although it is not normalized to be of unit length). As a map of the global action space, the frequency map is the Legendre transform associated to $`H`$ (cf. \[F.G, p.338\]).
(1.1.6) Definition We say that the simple surface of revolution $`(S^2,g)`$ is globally non-degenerate if $`\omega |_{\{H=1\}}`$ is an embedding.
This is the natural homogeneous analogue of the non-degeneracy condition of \[F.G, loc.cit\] to the effect that the Legendre transform be a global diffeomorphism, and has previously been studied in some detail by Bleher in the setting of simple surfaces of revolution \[Bl,§6\]. As will be seen in §1.3, the curve $`\{H=1\}`$ is the graph of a smooth function of the form $`I_2=F(I_1)`$ in the cone $`\mathrm{\Gamma }_o,`$ and the non-degeneracy condition (1.1.6) will follow as long as $`F`$ is a convex or concave function. In fact, in the proof of the Theorem we will only need to use that $`\{H=1\}`$ is non-degenerate at the one point $`(I_1,I_2)=(0,1).`$ This is sufficient because we assume the metric to be real analytic.
### 1.2 Length spectrum and periodic torii
We now come to the definition of length spectrum and simple length spectrum for a completely integrable geodesic flow. We first observe that the orbit thru $`(I,\varphi )`$ is periodic of period $`L`$ if and only if
$$L\omega _I=MZZ^2$$
$`(\mathrm{1.2.1}a)`$
for some $`M0.`$ The minimal positive such $`L`$ will be called the primitive period; the corresponding $`M`$ is known as the vector of winding numbers of the torus $`T_I`$. $`M`$ parametrizes the homology class of the closed orbit $`\gamma `$ since the latter has the form $`_{j=1}^2M_j\gamma _j(I)`$ relative to the homology basis $`\gamma _j(I)`$.
Due to the homogeneity, the period and vector of winding numbers are constant on the ray $`\mathrm{I}\mathrm{R}T_I`$. By Euler’s formula we then have
$$_IHI=\omega _II=H$$
hence the length is given in terms of the winding vector by
$$L=\frac{MI}{H(I)}$$
$`(\mathrm{1.2.1}b)`$
or simply $`L=MI`$ on the unit tangent bundle $`H=1`$.
It is clear that the periodicity condition $`L\omega =M`$ is independent of $`\varphi `$. Hence, all of the geodesics on $`T_I`$ are closed if any of them are. ( This also follows, of course, from the transitivity of the torus action on each invariant torus.) We therefore say:
(1.2.2) Definition
(a) A torus $`T_I`$ is a periodic torus if all the geodesics on it are closed.
(b) The period $`L`$ of the periodic torus is then the common period of its closed geodesics.
(c) The length spectrum $``$ of the completely integrable system is the set of these lengths.
(d) The completely integrable system has a simple length spectrum if there exist a unique periodic torus (up to time reversal) of each length $`L`$.
In the last statement (d) we are referring to the canonical involution $`\sigma :(x,\xi )(x,\xi )`$, which reverses the orientation of the geodesics. It is obvious that if $`T_I`$ is a periodic torus of period $`L`$, then so is $`\sigma (T_I)`$.
Let $`PerS_g^{}S^20`$ denote the set of periodic points for the geodesic flow on $`S_g^{}S^2`$, i.e. the set of points which lie on a closed geodesic. It is a union of periodic torii in $`S_g^{}S^2`$ together with points along the equatorial geodesics (which are degenerate torii). The set of all periodic points in $`T^{}S^20`$ is then equal to $`\mathrm{I}\mathrm{R}^+Per.`$ Since the invariant torii are parametrized by the points $`I\mathrm{\Gamma }`$ of action space, it is convenient to parametrize $`Per`$ by a subset of the level set $`\{H(I)=1\}\mathrm{\Gamma }.`$
(1.2.3) DefinitionThe set of points $`I\{H=1\}\mathrm{\Gamma }`$ such that $`T_IPer`$ will be called, with a slight abuse of notation, the set of periodic points on $`\{H=1\}`$ and will be denoted by $`𝒫`$. That is, $`𝒫=\{I\{H=1\}:L\mathrm{I}\mathrm{R}^+,L\omega _IZZ^2\}.`$
The following proposition will be needed later on (Proposition (4.1.4).
(1.2.4) Proposition If $`(S^2,g)`$ is non-degenerate (1.1.6), then $`𝒫`$ is dense in $`\{H=1\}`$.
Proof: Let $`𝒬:=\omega (𝒫)`$ equal the image of $`𝒫`$ under the frequency map $`\omega `$. Then by definition, $`𝒬`$ (for ‘rational points’) is the projection to the curve $`\omega (\{H=1\})`$ of the integer lattice in $`ZZ^2`$. It is clear that $`𝒬`$ is a dense set in $`\omega (\{H=1\})`$; since $`\omega `$ is an embedding, $`𝒫`$ is dense in $`\{H=1\}.`$
The next proposition shows that our inverse result is valid for a residual set of simple analytic surfaces of revolution.
(1.2.5) Proposition Let $`^{}`$ be the subset of metrics with simple length spectra. Then $`^{}`$ is a residual subset of $``$.
Proof: If $`L,L^{}`$, then there exist $`M,M^{}ZZ^2`$ and $`I,I^{}𝒬`$ such that $`L=MI,L^{}=M^{}I^{}.`$ So $`L=L^{}`$ implies that $`MIM^{}I^{}=0`$ hence that the $`I`$-coordinates of $`I,I^{}`$ are dependent over the rationals. Since the length spectrum moves continuously and non-trivially under deformations in $``$, such a dependence for fixed $`M,M^{}`$ can only hold on a closed nowhere dense set. The proposition follows.
### 1.3 First return times and angles
Let us consider more carefully the geometric interpretation of the Clairaut integral on a torus $`T_IS_g^{}S^2.`$ Since $`I_1`$ and $`H=|\xi |`$ are independent commuting coordinates, and since there are global action-angle variables, the different invariant torii in $`S_g^{}S^2`$ are parametrized by the values $`I_1=\iota [1,1]`$ of the Clairaut integral along $`\{H=1\}`$. Indeed by (1.1.3), the second action coordinate $`I_2`$ is determined from $`I_1`$ on $`\{H=1\}`$ by the formula
$$I_2=F(I_1):=|I_1|+\frac{1}{\pi }_{r_{}(I_1)}^{r_+(I_1)}\sqrt{1\frac{I_1^2}{a(r)^2}}𝑑r$$
$`(\mathrm{1.3.1})`$
where $`r_\pm (I_1)`$ are the roots of $`I_1^2=a(r)^2.`$ The torus in $`S_g^{}S^2`$ with $`I_1=\iota `$ is therefore $`T_{(\iota ,F(\iota ))}`$, which we will denote simply by $`T_\iota .`$ The projection of $`T_\iota `$ to $`S^2`$ is an annulus of the form $`r_+(\iota )<r<r_{}(\iota )`$. the geodesics on $`T_I`$ project to $`S^2`$ as almost periodic curves oscillating between the two extremal parallels $`r=r_+(\iota )`$ and $`r=r_{}(\iota ).`$
We observe then that $`T_\iota `$ contains a unique geodesic $`\gamma _\iota `$ which passes thru the intersection of the reference geodesics $`\gamma _M`$ and $`\gamma _E`$; $`\iota `$ equals $`a(r_o)cos\alpha `$ where $`\alpha (\iota )`$ is the angle between $`\gamma _\iota `$ and $`\gamma _M`$. In other words, $`\alpha (\iota )`$ is the common angle with which the geodesics on $`T_\iota `$ intersect the meridians as they pass thru the equator in the direction of the northern hemisphere. Since the length $`L_E`$ of the equator $`\gamma _E`$ equals $`2\pi \sqrt{a(r_o)}`$, and since this length is a symplectic invariant of the geodesic flow, the coordinates $`\iota `$ and $`\alpha `$ are related in an essentially universal fashion. Hence, either $`I_1`$ or $`\alpha (0,\pi )`$ could be used as an action coordinate on $`\{H=1\}`$; $`\alpha `$ is perhaps more geometric but $`I_1`$ is more convenient in calculations.
The picture is the same for any invariant torus $`T_I`$: Under the $`\mathrm{I}\mathrm{R}^+`$ action on $`T^{}S^20`$, it scales to a torus $`T_\iota S^{}S^2`$ and all features of its geometry are identical to that of $`T_\iota `$. Thus it carries a unique (parametrized) geodesic $`\gamma _I`$ such that $`\gamma _I(0)`$ is at the intersection $`\gamma _E\gamma _M`$, and so that $`\gamma _I^{}(0)`$ points to the northern hemisphere. The initial angle variables of $`(\gamma _I(0),\gamma _I^{}(0))`$ are therefore $`(\varphi _1(0),\varphi _2(0))=(0,0).`$ At time $`t`$ we denote the angle variables of $`(\gamma _I(t),\gamma _I^{}(t))`$ by $`(\varphi _1(t),\varphi _2(t))`$ where $`\varphi _1`$ essentially measures the meridianal angle and $`\varphi _2`$ measures the equatorial angle. We then introduce the following ‘first return times’:
(1.3.2) Definition
(Ei) The equatorial first return time is the minimal time $`\tau _E(I)>0`$ such that $`\varphi _1(\tau _E(I))=2\pi `$;
(Eii) The equatorial first return angle $`\omega _E(I):=\varphi _2(\tau _E(I))`$ is the change in angle along the equator of a geodesic on $`T_I`$ leaving $`\gamma _E\gamma _M`$ at $`t=0`$, upon its first return time to $`\gamma _E`$;
(Mi) The meridianal first return time is the minimal time $`\tau _M(I)`$ such that $`\varphi _2(\tau _M(I))=2\pi .`$
(Mii) The meridianal first return angle is the angle change $`\omega _M(I):=\varphi _1(\tau _M(I))`$ along the meridian of a geodesic on $`\gamma _M`$ leaving $`\gamma _E\gamma _M`$ at $`t=0`$, upon its first return time to $`\gamma _M`$.
The terminology ‘first return time’ is taken from dynamics. Note that $`\varphi _10(mod2\pi )`$ is the equation of the curve on $`T_I`$ which lies over the equator. Hence, $`\tau _E(I)`$ is the time of first return of $`\gamma _I(t)`$ to the equator in the direction of the northern hemisphere. This is actually the second time of intersection of $`\gamma _I`$ with the equator, the first one occurring when $`\gamma _I(T)`$ is heading to the southern hemisphere. This intersection is not in the projection of $`\varphi _10(mod2\pi )`$. Similarly, $`\varphi _20(mod2\pi )`$ is the equation of the fixed meridian $`\gamma _M`$, and so $`\tau _M`$ is the time of first return to the arc $`\gamma _M`$ (half of the closed geodesic).
The following gives some relations between the various angles and return times.
(1.3.3) Proposition Let $`\omega _I=(\omega _1,\omega _2)`$ be the frequency vector of the invariant torus $`T_I`$. Then:
(a) $`\tau _E\omega _2=\omega _E`$;
(b) $`\tau _M\omega _1=\omega _M`$;
(c) $`2\pi \frac{\omega _1(I)}{\omega _2(I)}=\omega _M,\mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}2}\pi \frac{\omega _2(I)}{\omega _1(I)}=\omega _E.`$
Proof: The equation of the geodesic $`(\gamma _I(t),\gamma _I^{}(t))`$ on $`T_I`$ beginning at $`(\varphi _1,\varphi _2)`$ in angle variables is $`(\varphi _1+t\omega _1(I),\varphi _2+t\omega _2(I))`$. When the torus is projected to the base, then $`\varphi _1`$ measures the meridianal angle and $`\varphi _2`$ measures the equatorial angle. So by definition of return times,
| $`(\varphi _1+\tau _E\omega _1,\varphi _2+\tau _E\omega _2)=(\varphi _1+2\pi ,\varphi _2+\omega _E)`$ |
| --- |
| $`(\varphi _1+\tau _M\omega _1,\varphi _2+\tau _M\omega _2)=(\varphi _1+\omega _M,\varphi _2+2\pi ).`$ |
The statements in the proposition follow immediately.
These first return times (and angles) are closely related to the (non-linear) Poincare maps of the geodesic flow. We recall that for each closed geodesic $`\gamma `$, the Poincare map $`𝒫_\gamma `$ is defined as the first return map of the geodesic flow, restricted to a symplectic transversal $`S_\gamma S_g^{}S^2`$ (surface of section). It is well-known $`𝒫_\gamma :S_\gamma S_\gamma `$ is a symplectic map \[K\]. Since the symplectic form on a cotangent bundle equals $`d\alpha `$ (with $`\alpha `$ the action form), it follows that $`𝒫_\gamma ^{}(\alpha )\alpha `$ is a closed 1-form on $`S_\gamma .`$ Since $`S_\gamma `$ may be assumed contractible, it follows that there exists a function $`\tau _\gamma `$ such that $`[𝒫_\gamma ^{}(\alpha )\alpha ]|_{S_\gamma }=d\tau _\gamma .`$ Since the integral of $`\alpha `$ over an arc of a unit speed geodesic just gives its length, $`\tau _\gamma `$ is the first return time of geodesics near $`\gamma `$ to $`S_\gamma .`$
In particular, let $`\gamma =\gamma _E`$ be the equator (in one of its orientations), and let $`S_E`$ denote a symplectic transversal at the point $`(\gamma _E(0),\gamma _E^{}(0))`$ thru the fixed meridian $`\gamma _M`$. Since $`\gamma _M`$ is transverse to the equator, we may define $`S_E`$ to consist of a small variation of $`\gamma _E^{}(0)`$ moved up and down a small arc of $`\gamma _M`$. We see then that the first return time $`\tau _{\gamma _E}`$ is precisely the first return time $`\tau _M`$ defined above. We also observe that the foliation of $`S^{}S^2`$ by invariant torii restricts to a foliation of $`S_E`$ by invariant circles for $`𝒫_{\gamma _E}`$, closing in on the center point where $`\gamma _E`$ intersects $`S_E`$. As noted above, the action coordinate $`I_1=\iota `$ gives a natural action (radial) coordinate on $`S_\gamma `$. Since $`S_\gamma S_g^{}S^2`$, the $`\mathrm{I}\mathrm{R}^+`$ homogeneity is broken and we may reformulate the twist condition (0.1 (iii)) as follows:
(1.3.4) Proposition The Poincare map $`𝒫_{\gamma _E}`$ is a twist map of $`S_E`$ iff $`\omega _M^{}0`$, where $`\omega _M^{}=\frac{}{\iota }\omega _M.`$
Proof: The coordinates $`(I_1=\iota ,\varphi _1)`$ restrict to a system of symplectic coordinates on $`S_E`$, in terms of which the Poincare map takes the form
$$𝒫_{\gamma _E}(\iota ,\varphi _1)=(\iota ,\varphi _1+\tau _M\omega _1)=(\iota ,\varphi _1+\omega _M).$$
By definition it is a twist map if $`\omega _M^{}0`$ in a neighborhood of $`\iota =0`$ in $`S_E`$.
For background on the twist condition in a related context see \[F.G\]\[P\].
The situation for the other periodic orbits is different since they come in one-parameter families. Thus, for the (closed) meridian geodesic $`\gamma _M`$, a transversal $`S_M`$ is given by the equator $`\gamma _E`$ and a small variation of $`\gamma _M^{}(0)`$ along it. The foliation by invariant torii restricts to $`S_M`$ to a foliation by invariant lines (non-closed curves), including the curve of closed geodesics thru $`\gamma _M`$. The Poincare map $`P_{\gamma _M}`$ is then of parabolic type; the first return time is $`\tau _E`$ above.
For the inverse problem it will be necessary to have expressions for these return times and angles in terms of the metric. This will also make the twist condition more transparent. For ease of quotation, it is convenient to make a change of dependent and independent variables, following \[Besse\] and Darboux \[D\]. Equivalent expressions in the original polar coordinates can also be easily derived, and will be given below.
(1.3.5) Proposition Suppose that $`(S^2,g)`$ is a simple surface of revolution. Then there exists a coordinate system $`(u,\theta )`$ on $`U`$, with $`sinu=a(r)`$ and a smooth function $`f`$ on $`[1,1]`$ such that $`f(1)=1,f(1)=1`$ and such that
$$g=[f(cosu)]^2du^2+sin^2ud\theta ^2.$$
Proof: First, define
$$b(r)=\{\begin{array}{cc}sin^1(a(r))\hfill & r[0,r_o]\hfill \\ \pi sin^1(a(r))\hfill & r[r_o,L]\hfill \end{array}$$
and
$$c(v)=\{\begin{array}{cc}(a|_{[0,r_o]})^1(\sqrt{1v^2})\hfill & v[0,1]\hfill \\ (a|_{[r_o,L]})^1(\sqrt{1v^2})\hfill & v[1,0]\hfill \end{array}$$
Then define $`f:(1,1)\mathrm{I}\mathrm{R}`$ by
$$\{\begin{array}{c}f(v)=\frac{v}{a^{}[c(v)]},(v0)\hfill \\ f(0)=\frac{1}{a^{^{\prime \prime }}(r_o)}\hfill \end{array}$$
Since $`b(r)=u,c(cosu)=b^1(u)=r`$ we have $`a(r)=sinu,a^{}[c(cosu)]dr=cosudu`$ and $`b^{}(r)dr=du.`$ The smoothness properties of $`f`$ follow from those of $`a`$ and from the fact that $`a(r)`$ and $`sinu`$ have the same qualitative shapes \[Besse, loc.cit.\].
The geometric result is the following \[Besse, Theorem 4.11\]:
(1.3.6) Proposition In the above coordinates, the equator $`\gamma _E(s)=(u(s),\theta (s))`$ has the equation
$$u(s)\frac{\pi }{2},\theta (s)=s.$$
Any other geodesic $`\gamma (s)=(u(s),\theta (s))`$ in $`U`$ is contained between two parallels $`u=i`$ and $`u=\pi i`$ and the angle $`\theta (i)`$ between two consecutive points of contact with these parallels is given by:
$$\theta (i)=sini_i^{\pi i}\frac{f(cosu)}{sinu(sin^2usin^2i)^{\frac{1}{2}}}𝑑u.$$
The length of this arc of $`\gamma `$ is given by
$$s(i)=_i^{\pi i}\frac{sin(u)f(cosu)}{(sin^2usin^2i)^{\frac{1}{2}}}𝑑u.$$
Sketch of Proof: Using the Clairaut integral, the equations of the geodesic have the form:
$$\{\begin{array}{c}\frac{d\theta }{ds}=\frac{sini}{sin^2u}\hfill \\ \frac{d\theta }{du}=\frac{sin(i)f(cosu)}{sinu(sin^2usin^2i)^{\frac{1}{2}}}\hfill \\ \frac{ds}{du}=\frac{sin(u)f(cosu)}{(sin^2usin^2i)^{\frac{1}{2}}}\hfill \end{array}$$
The formulae above follow by integration.
The following is geometrically obvious:
(1.3.7) Proposition Let $`T_IS_g^{}S^2`$ denote an invariant torus with $`H(I)=1`$ and let $`i(I)`$ be the $`u`$-coordinate of the extremal parallel closest to $`N`$ in the projection of $`T_I`$ to $`S^2`$. Then:
(i) $`\tau _E(I)=2s(i(I))`$;
(ii) $`\omega _E(I)=2(\theta (i(I))1)`$;
(iii) $`\tau _E(I)=\frac{1}{\pi }_{r_{}(I_1)}^{r_+(I_1)}(1\frac{I_1^2}{a(r)^2})^{\frac{1}{2}}𝑑r`$
(iv) $`\omega _E(I)=1+\frac{I_1}{\pi }_{r_{}(I_1)}^{r_+(I_1)}a(r)^2(1\frac{I_1^2}{a(r)^2})^{\frac{1}{2}}𝑑r.`$
Proof:
(i) By definition, $`\omega _E(I)`$ is the change in angle along the equator between a geodesic $`\gamma _I`$ on $`T_I`$, starting at the equator on a fixed meridian and heading towards the northern hemisphere, and the fixed meridian, upon second intersection with the equator. We can view this arc of the geodesic as consisting of three pieces: one from the equator to the northern extremal parallel, one on the ‘back-side’ between the two extremal parallels, and one on the ‘front-side’ from the southern parallel to the equator. Since the lengths of the ‘front-side’ arcs are unchanged by rotation, we can rotate one until the two make up a smooth geodesic arc between the parallels. The length of the geodesic arc is therefore twice that of an arc between the parallels, i.e. $`\tau _E(I)=2s(i(I)).`$
(ii) For the angle change: In the same way, the change in $`\theta `$ along this arc of $`\gamma _I`$ is the change in $`\theta `$ of two arcs between the extremal parallels. We subtract $`1`$ since $`\omega _I`$ measures the addition to one full revolution.
(iii) Since $`sin(u)=a(r)`$ and $`f(cosu)du=dr`$, we have
$$2s(i(I))=\pi I_1_{i(I)}^{\pi i(I)}\frac{dr}{\sqrt{1\frac{I_1^2}{a(r)^2}}}.$$
(iv) From Propositions 1.3.2 we have that $`\frac{\omega _1}{\omega _2}=\frac{1}{2\pi }\omega _E.`$ Since $`H(I_1,I_2)=1`$ implies that $`\omega _1+F^{}(I_1)\omega _2=0,`$ we get from (1.3.1) that $`\omega _E=2\pi F^{}(I_1).`$
(1.3.8) Corollary The non-degeneracy condition (1.1.6) is satisfied if $`𝒫_{\gamma _E}`$ is globally twisted, i.e. if $`\omega _M^{}>0`$ or $`\omega _M^{}<0.`$.
Proof: In these cases, $`F`$ is convex (or concave). Since the set $`\{H=1\}`$ is the graph of $`F`$ in $`\mathrm{\Gamma }_o`$, the Gauss map (and hence the frequency map) is an embedding.
Remark Both cases, of concavity and convexity, occur for ellipsoids of revolution, see \[Bl\]. The separating case of the round sphere is of course degenerate.
### 1.4 Classical Birkhoff invariants
The classical Birkhoff normal form of Hamiltonian $`H`$ near a non-degenerate periodic orbit $`\gamma `$ is a germ of a completely integrable system to which $`H`$ is symplectically equivalent in a ‘formal neighborhood’ of $`\gamma `$ (see e.g. \[F.G\] for a detailed discussion). In the case at hand, where $`H`$ is already completely integrable, the Birkhoff normal form is simply $`H`$ itself expressed in terms of action-angle variables. The ‘Birkhoff invariants’ of the Hamiltonian at a torus $`T_{I_o}`$ may then be identified with the germ of $`H(I_1,I_2)`$ or of $`\omega _I`$ at $`I=I_o.`$
Since $`H`$ is homogeneous of degree 1, it is equivalent and somewhat clearer to define the Birkhoff invariants after first breaking the $`\mathrm{I}\mathrm{R}^+`$ symmetry. That is, we would like to introduce a ‘base’ to the cone $`\mathrm{\Gamma }_o`$. The most natural one may appear to be the energy level $`\{H=1\}`$; but for the purpose of calculating wave invariants at $`T_{I_o}`$ it is more convenient to use the tangent line $`\omega _{I_o}(II_o)=0`$ at a point $`I_o\{H=1\}.`$
Let us first consider the level set $`\{H=1\}`$ as the transversal. The homogeneous function $`H`$ is obviously determined by the curve $`H(I_1,I_2)=1`$ whose equation is given by (1.3.1). Hence we can define the Birkhoff invariants at a torus $`T_{I_o}`$ to be the Taylor coefficients of the function $`F^{}`$ at $`I_o\{H=1\}`$. By (1.3.7 (iv)) it is equivalent to put:
(1.4.1) Definition The (first) Birkhoff invariants of $`H`$ at an invariant torus $`T_\iota `$ with $`\iota \{H=1\}`$ are the Taylor coefficients of $`\omega _E`$ at $`\iota `$ in the coordinate $`I_1`$.
Secondly, let us consider the tangent lines as transversals: We fix a point $`I^o\{H=1\}`$, let $`\omega ^o`$ denote the common frequency vector of the ray of torii $`\mathrm{I}\mathrm{R}^+T_{I^o}`$ and put
$$I\omega ^o:=\underset{k=1}{\overset{2}{}}\omega _k^oI_k.$$
$`(\mathrm{1.4.2})`$
The equation of the tangent line to $`\{H=1\}`$ at $`I^o`$ in the action cone $`\mathrm{\Gamma }_o`$ is then given by $`\omega ^oI=1`$. Note that $`I\omega ^o`$ is homogeneous of degree 1 and hence equals $`H(I)`$ along the ray $`\mathrm{I}\mathrm{R}I^o`$; consequently it is elliptic (non-vanishing) in a conic neighborhood $`W_o\mathrm{\Gamma }_o`$ of it. The conic neighborhood will be parametrized in the following way: we fix a basis (i.e. a non-zero vector) $`v`$ of the line $`I\omega ^o=0`$, and define the map
$$(\rho ,\xi )\rho (I^o+\xi v),\xi (ϵ,ϵ).$$
$`(\mathrm{1.4.3}a)`$
For sufficiently small $`ϵ`$, this map sweeps out a conic neighborhood $`W_o`$ of $`I^o`$ with inverse given by
$$\rho =\omega ^oI,\xi v_j:=\frac{I_j}{I\omega ^o}I_j^o.$$
$`(\mathrm{1.4.3}b)`$
Since $`H(I_1,I_2)=(\omega ^oI)H(\frac{I_1}{\omega ^oI},\frac{I_2}{\omega ^oI})`$ and since $`(\frac{I_1}{\omega ^oI},\frac{I_2}{\omega ^oI})\{\omega ^oI=1\}`$ we may write
$$H(I_1,I_2)=\rho h_{I^o}(\xi )$$
$`(\mathrm{1.4.4}a)`$
where $`h_{I^o}`$ is the function on $`W_o\{\omega ^oI=1\}`$ defined by
$$h_{I^o}(\xi ):=H(I^o+\xi v).$$
$`(\mathrm{1.4.4}b)`$
The $`C^{\mathrm{}}`$ Taylor expansion of $`h_{I^o}(\xi )`$ around $`\xi =0`$ is then a symplectic invariant of $`H`$.
(1.4.5) Definition The second (tangential) classical Birkhoff invariants of $`H`$ associated to the periodic torus $`T_{I^o}`$ are the Taylor coefficients $`h_{I^o}^\alpha (0).`$
In the real analytic case the Taylor coefficients determine $`h_{I^o}`$ and hence $`H`$ by homogeneity. It is more or less obvious that the first and second Birkhoff invariants also carry the same information in the $`C^{\mathrm{}}`$ case. To be sure, we prove:
(1.4.6) Proposition The first Birkhoff invariants canonically determine the tangential Birkhoff invariants and vice versa.
Proof:
By definition
$$h_{I^o}(\xi )=H(I^o+\xi v).$$
Therefore,
$$H(h_{I^o}(\xi )^1(I^o+\xi v))=1$$
or equivalently
$$h_{I^o}(\xi )^1(I_2^o+\xi v_2)=F(h_{I^o}(\xi )^1(\xi v_1+I_1^o)).$$
Writing $`u=\frac{I_1^o+\xi v_1}{h_{I^o}}`$, this says
$$\frac{F(u)}{u}=\frac{I_2^o+\xi v_2}{I_1^o+\xi v_1}.$$
Hence the knowledge of Taylor coefficients of $`F`$ is equivalent to knowledge of the Taylor coefficients of $`h_{I^o}`$.
We may reformulate the non-degeneracy condition (0.1 (iii)) in terms of $`h_{I^o}`$ where $`T_{I^o}`$ is the meridian torus:
(1.4.7) Proposition $`(S^2,g)`$ satisfies the non-degeneracy condition (0.1 (iii)) as long as $`\alpha :=h_{I^o}^{\prime \prime }(0)0.`$
Proof: By definition, $`h_{I^o}(\xi )=H(I^o+\xi v)`$. At the meridian torus, $`v=(1,0)`$ so $`h_{I^o}(\xi )=H(I^o+(\xi ,0))`$ and $`h_{I^o}^{}(\xi )=\frac{}{I_1}H(I^o+(\xi ,0))=\omega _1(I^o+(\xi ,0)`$. Hence
$$h_{I^o}^{\prime \prime }(0)=\frac{}{I_1}\omega _1(I^o).$$
Since $`\omega _1(I^o)=0`$ it follows that
$$\frac{}{I_1}\omega _M(I^o)=\frac{\omega _1^{}(I^o)}{\omega _2(I^o)}.$$
Hence, the condition $`h_{I^o}^{\prime \prime }(0)0`$ is equivalent to the condition that $`\frac{}{\iota }\omega _M(I^o)0.`$ The equivalence of this to (0.1 (iii)) is proved in Propositions (1.3.4) and again in Corollary (1.3.8).
### 1.5 Symplectic conjugacy of geodesic flows
The Birkhoff normal form of a Hamiltonian $`H`$ at a closed orbit (or family of closed orbits) is a symplectic conjugacy invariant of $`H`$ in a neighborhood of the orbit(s). Hence the global Birkhoff normal form $`H(I_1,I_2)`$ of a completely integrable Hamiltonian is a symplectic conjugacy invariant. The purpose of this section is to show that it is a complete conjugacy invariant. We begin by showing that the homogeneous Hamiltonian torus actions commuting with geodesic flows of simple surfaces of revolution are all symplectically equivalent:
(1.5.1) Proposition Suppose $`(S^2,g_1)`$ and $`(S^2,g_2)`$ are smooth surfaces of revolution of simple type and let $`(I_1,I_2)`$ resp. $`(J_1,J_2)`$ denote their global action variables as above. Then there exists a homogeneous symplectic diffeomorphism $`\chi :T^{}S^2T^{}S^2`$ such that $`\chi ^{}J_i=I_i`$.
Sketch of Proof : Let $`\varphi _i`$, resp. $`\psi _i`$ be the dual angle variables on $`(S^2,g_1)`$ resp. $`(S^2,g_2)`$. Except on the equators of $`(S^2,g_i)`$ a point in $`T^{}S^2`$ is uniquely specified by its action-angle coordinates. Define $`\chi `$ to be the identity map in action-angle coordinates with respect to the two metrics. It is obvious that $`\chi `$ is a homogeneous symplectic diffeomorphism on the complement of the equators, so to prove the Corollary it suffices to show that $`\chi `$ extends to the equators with this property.
Since $`\chi `$ is homogeneous of degree 1, it is necessary and sufficient to define it on the unit cotangent bundles. Moreover, since it commutes with the Hamilton flow of $`p_\theta `$ on the regular set, its extension must also do so. Hence it must be the lift of a diffeomorphism $`\overline{\chi }`$ on the orbit space $`𝒪:=S^{}S^2/S^1`$ of the rotation action. This action is free so the natural projection $`p:S_g^{}S^2𝒪`$ must be diffeomorphic to the standard projection from $`S0(3)S^2.`$ The image of the torus foliation defined by level sets of $`p_\theta `$ is a singular foliation of $`𝒪`$ formed by level sets of the function $`\overline{p}_\theta `$ induced by $`p_\theta `$ on $`𝒪`$, and the two singular points $`o^\pm `$ are the images of the equators $`\gamma _E^\pm `$. Since $`\overline{p}_\theta `$ is, by assumption, a perfect Morse function on $`𝒪S^2`$ for each metric, the quotient map $`\overline{\chi }`$ on the punctured quotient $`𝒪\{o^\pm \}`$ of $`\overline{\chi }`$ extends smoothly to the completion. It follows that $`\chi `$ extends smoothly as a rotationally equivariant map on the completion of $`S^{}S^2\{\gamma _E^\pm \}`$ and the homogeneous extension must be symplectic. (See \[CV.1\] for more on the behaviour near the poles).
Remark: In fact, this proposition can be sharpened to say: there exists only one homogeneous Hamiltonian $`S^1\times S^1`$ action on $`T^{}S^20`$ (up to symplectic equivalence). This follows from a homogeneous analogue of the Delzant classification of completely integrable torus action on compact Kahler manifolds. We hope to report on more general results of this kind at a later time.
Thus, the torus actions defined above on the cotangent bundles of simple surfaces of revolution are always symplectically equivalent. The question arises when the geodesic flows are symplectically equivalent. The answer can be given by expressing the norm functions $`|\xi |_g`$ of the metrics in terms of the global action variables. Before doing so, we note that the action variables are not quite uniquely defined above because the choice of generators $`\gamma _i(b)`$ is not unique. For instance, one might have permuted the roles of $`N`$ and $`S`$. Hence we have:
(1.5.2) Proposition Let $`(S^2,g_i)`$ be a simple surfaces of revolution, and let $`(I_1,I_2)`$ resp. $`(J_1,J_2)`$ be the global action variables defined above. Let $`|\xi |_{g_1}=H_1(I_1,I_2)`$ resp. $`|\xi |_{g_2}=H_2(J_1,J_2)`$ be the expressions of the metric norms of $`g_i`$ in terms of action variables. Then: the geodesic flows of $`(S^2,g_i)`$ are homogeneously symplectically equivalent if and only if there a linear map $`A=\left(\begin{array}{cc}a_{11}\hfill & a_{12}\hfill \\ a_{21}\hfill & a_{22}\hfill \end{array}\right)SL(2,ZZ)`$ such that $`H_1=H_2A.`$
Proof If such a choice exists, then the map $`\chi `$ above obviously defines a symplectic conjugacy.
Conversely, suppose the geodesic flows are symplectically conjugate by a homogeneous symplectic diffeomorphism $`\chi `$, i.e. $`\chi ^{}H_2(J_1,J_2)=H_1(I_1,I_2).`$ Then $`\chi ^{}J_i`$ are global action variables for the geodesic flow of $`(S^2,g_1)`$. But global action variables are almost unique: they correspond to a trivialization of the lattice bundle $`H_1(F_b,ZZ)`$ \[D\]\[G.S\]. Therefore there exists $`ASL(2,ZZ)`$ so that $`\chi ^{}J=AI.`$ Then $`H_1(I)=H_2(AI)`$.
The following gives a geometric interpretation of the conjugacy condition:
(1.5.3) Proposition Suppose $`g_1,g_2`$. Then their geodesic flows are symplectically equivalent if and only if their equatorial first return times $`\tau _E`$ and angles $`\omega _E`$ are equal.
Proof: Suppose first that the flows are symplectically equivalent. After choosing compatible bases for the homology we may then assume that the expressions for $`H_1`$ and $`H_2`$ in global action-angle variables are the same. Then the frequency maps are the same and hence by Proposition (1.3.3) the equatorial return angles are the same. Also, the Poincare maps $`𝒫_{\gamma _E}`$ are conjugate and hence the equatorial return times are equal.
Conversely, if the equatorial return times and angles are the same, then the flows have the same frequency vectors (as functions of the global action angle variables) and hence have the same global Birkhoff normal forms. By Proposition (1.5.2) the flows are symplectically equivalent.
## 2 Quantum dynamics and normal form
It is owing to the following notion that simple surfaces of revolution are so manageable.
(2.1) Definition The wave group $`e^{it\sqrt{\mathrm{\Delta }}}`$ of a compact, Riemannian n-manifold $`(M,g)`$ is quantum torus integrable if there exists a unitary Fourier-Integral representation
$$\widehat{\tau }:T^nU(L^2(M)),\widehat{\tau }_{(t_1,\mathrm{},t_n)}=e^{i(t_1\widehat{I_1}+\mathrm{}t_n\widehat{I_n})}$$
of the n-torus and a symbol $`\widehat{H}S^1(\mathrm{I}\mathrm{R}^n0)`$ such that $`\sqrt{\mathrm{\Delta }}=\widehat{H}(I_1,\mathrm{},I_n).`$
In the formula above, we follow physics notation in indicating operators (as opposed to their symbols) with a ‘hat’. Thus, the generators $`\widehat{I}_j`$ are first order pseudodifferential operators with the property that $`e^{2\pi i\widehat{I}_j}=C_jId`$ for some constant $`C_j`$ of modulus one.
Since $`\widehat{H}`$ is a first order elliptic symbol on $`\mathrm{I}\mathrm{R}^n0`$ it has an asymptotic expansion in homogeneous functions of the form:
$$\widehat{H}H_1+H_o+H_1+\mathrm{},H_j(rI)=r^jH_j(I).$$
$`(2.2)`$
The quantum action operators are uniquely defined up to a choice of basis of $`H^1(\mathrm{I}\mathrm{R}^n/ZZ^n,ZZ)`$, the terms $`H_j`$ are uniquely determined up to the same ambiguity.
The principal symbols $`I_j`$ of the $`\widehat{I}_j`$’s generate a classical Hamiltonian torus action, so any quantum torus action automatically induces a classical one. Conversely, it is a theorem of Boutet de Monvel-Guillemin and Weinstein that any classical Hamiltonian torus action can be quantized \[BM.G, Appendix, Proposition 6.6\]. Since metrics in $``$ give rise to quantum torus actions, we have, in particular:
(2.3) Proposition (cf. \[CV.1\]) Suppose $`g`$. Then $`\sqrt{\mathrm{\Delta }_g}`$ is quantum torus integrable.
We also observe that the following holds for any Laplacian commuting with a quantum torus action:
(2.4) Proposition For any $`\mathrm{\Delta }`$, $`H_o=0.`$
Proof: Since the subprincipal symbol $`\sigma _{sub}(\sqrt{\mathrm{\Delta }})`$ equals zero, the same is true $`\widehat{H}:=\widehat{H}(\widehat{I}_1,\widehat{I}_1).`$ Now $`\sigma _{sub}(\widehat{H})`$ is invariantly defined (independent of a choice of symplectic coordinates); hence it may be expressed in action-angle coordinates in the form
$$0=\sigma _{sub}(\widehat{H})=H_o\frac{1}{i}\underset{j}{}\frac{^2H_1}{I_j\varphi _j}.$$
But the mixed derivative term automatically vanishes since $`H_1`$ is a function only of the classical action variables.
Let us specialize to the case of $`\sqrt{\mathrm{\Delta }_g}`$ with $`g.`$ From the fact that $`e^{2\pi i\widehat{I}_j}=C_jId`$ for a quantum torus action, it follows that the joint spectrum of the quantum moment map
$$Sp()ZZ^2\mathrm{\Gamma }+\{\mu \}$$
is the set of integral lattice points, translated by $`\mu `$, in the closed convex conic subset $`\mathrm{\Gamma }\mathrm{I}\mathrm{R}^2`$. The vector $`\mu =(\mu _1,\mu _2)`$ can be identified with the Maslov indices of the homology basis of the invariant torii. In the case of $`\sqrt{\mathrm{\Delta }_g}`$, $`\mu =(0,1/2)`$ \[CV.1\].
Expressing $`\sqrt{\mathrm{\Delta }_g}`$ in the form $`\widehat{H}(\widehat{I}_1,\widehat{I}_2)`$ we have that
$$Sp(\sqrt{\mathrm{\Delta }_g})=\{\widehat{H}(N+\mu ):NZZ^2\mathrm{\Gamma }_o\}.$$
Thus, the eigenvalues of $`\sqrt{\mathrm{\Delta }_g}`$ have the form:
$$\lambda _NH_1(N+\mu )+H_1(N+\mu )+\mathrm{}$$
and the wave trace takes the form
$$Tre^{it\sqrt{\mathrm{\Delta }_g}}=\underset{NZZ^2}{}e^{it\widehat{H}(N+\mu )}.$$
$`(2.5)`$
The symbol $`\widehat{H}`$ may be regarded as a global quantum Birkhoff normal form. As in the classical case, it has a germ at any periodic torus, so these may be regarded as the microlocal Birkhoff canonical forms. To be more precise, we imitate the definition of the classical tangential Birhoff normal forms and write
$$H_j(I_1,I_2)=(\omega ^oI)^jH_j(\frac{I_1}{\omega ^oI},\frac{I_2}{\omega ^oI}):=(\omega ^oI)^jh_j(\xi )$$
$`(2.6)`$
where $`\xi `$ is a linear coordinate relative to a vector $`v`$ generating $`\omega ^oI=1.`$ We then Taylor expand $`h_j(\xi )`$ around $`\xi =0`$:
$$h_j(\xi )=\underset{\alpha 0}{}h_j^\alpha (0)\xi ^\alpha .$$
(2.7) Definition The homogeneous quantum Birkhoff normal form coefficients of $`\widehat{H}`$ at the periodic torus $`T_{I^o}`$ are the Taylor coefficients $`h_j^\alpha `$ for $`j=1,0,1,\mathrm{}.`$
(2.8) Proposition Assume that $`g^{}`$. Then all of the functions $`H_j`$ are real analytic in the interior of the cone $`\mathrm{\Gamma }_o`$ and all of the functions $`h_j`$ are real analytic near $`\xi =0`$.
Proof: First consider $`|\xi |_g=H_1(I_1,I_2)=H(I_1,I_2).`$ We know that $`H`$ is a $`C^{\mathrm{}}`$ homogeneous function on $`\mathrm{\Gamma }_o`$. On the other hand, from the formula $`I_2=G(|\xi |_g,I_1)`$ we see that $`I_2`$ is a real analytic function of $`|\xi _g|,I_1.`$ Since $`G`$ is the inverse function to $`H`$ with respect to the first variable, $`H`$ must also be a real analytic function of $`I_1,I_2.`$
Next recall that $`\widehat{I}_2`$ is defined (up to a smoothing term) as the function $`\widehat{G}(\sqrt{\mathrm{\Delta }},\widehat{I}_1)`$ which has principal symbol $`G(|\xi |_g,I_1)`$ and which has integral spectrum. More precisely, one begins with $`G(\sqrt{\mathrm{\Delta }},\widehat{I}_1)`$, which has the property that $`exp(2\pi iG(\sqrt{\mathrm{\Delta }},\widehat{I}_1)+\mu _j)=I+K`$ with $`K`$ of order -1. One defines $`R=\frac{1}{2\pi i}Log(I+K)`$ and puts $`\widehat{G}(\sqrt{\mathrm{\Delta }},\widehat{I}_1)=G(\sqrt{\mathrm{\Delta }},\widehat{I}_1)+\mu _j+R_j.`$
Since $`G`$ is a real analytic function and we only apply the holomorphic functional calculus in the steps of the construction, it follows that $`\widehat{G}`$ is an analytic function of $`(\sqrt{\mathrm{\Delta }},\widehat{I}_1)`$.
Since $`\widehat{I}_2=\widehat{G}(\sqrt{\mathrm{\Delta }},\widehat{I}_1)`$ has the inverse function $`\sqrt{\mathrm{\Delta }}=\widehat{H}(\widehat{I}_1,\widehat{I}_2)`$, it follows again by the inverse function theorem for analytic functions that $`\widehat{H}`$ is a real analytic function.
The real analyticity of the $`h_j`$’s follows from that of the $`H_j`$’s.
## 3 Wave invariants as non-commutative residues
To relate the wave invariants to the coefficient of the normal form, it will be helpful (as in \[G.1\]\[Z.1,2\]) to express the wave invariants as non-commuative residues of the wave operator and its time-derivatives. Let us recollect how this goes.
The non-commutative residue of a Fourier Integral operator is an extension of the well-known non-commutative residue of a pseudodifferential operator $`A`$ on a compact manifold $`M`$ , defined by
$$\text{res}(A)=2\text{Res}_{s=0}\zeta (s,A)$$
where
$$\zeta (s,A)=\text{Tr}A\mathrm{\Delta }^{s/2}(\text{Re }s>>0).$$
Here, $`\mathrm{\Delta }`$ is a Laplacian (or any positive elliptic operator) on $`M`$. From work of Seeley, Wodzicki and Guillemin, one knows that $`\zeta (s,A)`$ is holomorphic in Re $`s>\frac{1}{2}dimM+`$ ord$`(A)`$ and admits a meromorphic continuation to $`\mathrm{C}`$, with simple poles at $`s=dimM+`$ord $`(A)k`$ ($`k=0,1,2,\mathrm{}`$). The residue at $`s=0`$ has a number of remarkable properties (not shared by the residues of the other poles):
| | | res$`(AB)=`$ res$`(BA)`$, i.e. res is a trace on the algebra $`\mathrm{\Psi }^{}(M)`$ of pseudodifferential operators over $`M`$; |
| --- | --- | --- |
| | | res$`(A)`$ is independent of the choice of $`\sqrt{\mathrm{\Delta }}`$; |
| | | there is a local formula for the residue, res$`(A)=(2\pi )^n_{S^{}M}a_n(x,\xi )i_{}𝑑xd\xi `$, |
where $`a_n`$ is the term of degree $`(n)`$ in the complete symbol expansion $`a{\displaystyle \underset{m}{\overset{\mathrm{}}{}}}a_j`$ for $`A`$; $`dxd\xi `$ is the canonical symplectic volume measure on $`T^{}M`$.
These results may be extended to Fourier integral operators as follows: Let $`A`$ be a Fourier Integral operator in $`I^m(M\times M,\mathrm{\Lambda })`$ for some homogeneous canonical relation $`\mathrm{\Lambda }T^{}(M\times M)\backslash 0`$ and $`mZZ`$. Below, diag$`(X\times X)`$ denotes the diagonal in $`X\times X`$. Below we will sketch a proof of the following:
(3.1) Theorem Suppose $`\mathrm{\Lambda }`$ and diag$`(T^{}M\times T^{}M)`$ intersect cleanly. Then $`\zeta (s,A)=`$ Tr $`A\mathrm{\Delta }^{s/2}(\text{Res }>>0)`$ has a meromorphic continuation to $`\mathrm{C}`$, with simple poles at $`s=m+1+\frac{e_01}{2}j`$, where $`e_0=dim\mathrm{\Lambda }`$ diag$`(S^{}M\times S^{}M)`$, and $`j=0,1,2,\mathrm{}`$ .
The clean intersection hypothesis above is that $`\mathrm{\Lambda }`$ diag$`(T^{}M\times T^{}M)\backslash 0`$ is a clean intersection. It is satisfied in the case where $`\mathrm{\Lambda }=C_t`$ if and only if the fixed point set of $`G^t`$ is clean. Hence, Theorem (3.1) implies:
(3.2) Corollary Let $`\zeta (s,t)=`$ Tr $`U(t)\mathrm{\Delta }^{s/2}`$.
If the fixed point set of $`G^t`$ is clean, then $`\zeta (,t)`$ has a meromorphic continuation to $`\mathrm{C}`$, with simple poles at $`s=1+\frac{dimS\text{ Fix}(G^t)1}{2}j`$ ($`j=0,1,\mathrm{},2`$).
Here, $`S`$ Fix$`(G^t)`$ is the set of unit vectors in the fixed point set of $`G^t`$. In the case of a completely integrable system, $`S`$ Fix$`(G^L)`$ is the union of the periodic torii with period $`L`$. We assume here, and henceforth, that the periodic torii are all clean fixed point sets for $`G^t`$ on $`S^{}M`$.
The non-commutative residue of the Fourier integral operator $`A`$ is then defined by:
(3.3) Definition
$$\text{res}(A):=\text{Res}_{s=0}\zeta (s,A)$$
just as in the case of pseudodifferential operators. And just as in that case, res$`(A)`$ has some remarkable properties:
\- it is independent of the choice of $`\mathrm{\Delta }`$;
\- if either $`A`$ or $`B`$ is associated to a local canonical graph, then res$`(AB)=`$ res$`(BA)`$
\- there is a local formula for res$`(A)`$.
The basic properties of res$`(A)`$ may be deduced from a singularity analysis of the closely related distribution trace $`S(t,A):=`$ Tr$`AU(t)`$ (cf. \[Z.4\]). Under the cleanliness assumption above, $`S(t,A)`$ is a Lagrangean distribution on $`\mathrm{I}\mathrm{R}`$ with singularities at the set of ‘sojourn times’,
$$𝒮𝒯=\{T:(x,\xi )S^{}M:(x,\xi ,G^T(x,\xi ))\mathrm{\Lambda }\}.$$
For a given sojourn time $`T`$, the corresponding set $`W_T=\{(x,\xi ):(x,\xi ,G^t(x,\xi ))\mathrm{\Lambda }\}`$ of sojourn rays fills out a submanifold of $`S^{}M`$ of some dimension $`e_T`$. With $`N(\lambda ,A)=_{\sqrt{\lambda _j}\lambda }(A\phi _j,\phi _j)`$, we have
(3.4) Proposition \[Z.4, Proposition 1.10\]. If $`\rho C^{\mathrm{}}(\mathrm{I}\mathrm{R})`$ with $`\widehat{\rho }C_0^{\mathrm{}}(\mathrm{I}\mathrm{R})`$, supp $`(\widehat{\rho })𝒮𝒯=\{0\}`$ and $`\widehat{\rho }1`$ near 0, then
$$\rho dN(\lambda ,A)C_n\lambda ^{m+\frac{e_01}{2}}\alpha _j\lambda ^j,$$
where $`C_n`$ is a universal constant. The coefficients have the form,
$$\alpha _j=_{\mathrm{\Lambda }_\mathrm{\Delta }}\omega _j𝑑\lambda _\mathrm{\Delta }$$
where $`\mathrm{\Lambda }_\mathrm{\Delta }=\mathrm{\Lambda }`$ diag$`(S^{}M\times S^{}M)`$, $`e_0=dim\mathrm{\Lambda }_\mathrm{\Delta }`$, $`d\lambda _\mathrm{\Delta }`$ is a canonical density on $`\mathrm{\Lambda }_\mathrm{\Delta }`$ and the functions $`\omega _j`$ are determined by the $`j`$-jet of the (local) complete symbols of $`AU`$ along $`\mathrm{\Lambda }_\mathrm{\Delta }`$.
Proof of Theorem (3.1). As in the pseudodifferential case \[DG, Proposition 2.1\], we have
$$\begin{array}{ccc}\text{Tr}A\mathrm{\Delta }^{s/2}\hfill & =\hfill & \chi _s,dN(^\text{.},A)\hfill \\ & =\hfill & \chi _s,\rho dN(^\text{.},A)+\text{ entire}.\hfill \end{array}$$
Hence,
$$\begin{array}{ccc}\zeta (s,A)\hfill & =\hfill & C_n\underset{j=0}{\overset{\mathrm{}}{}}\alpha _j_1^{\mathrm{}}\lambda ^{\frac{e_01}{2}+msj}𝑑\lambda +\text{ entire}\hfill \\ & =\hfill & C_n\underset{j=0}{\overset{\mathrm{}}{}}\frac{\alpha _j}{m+1+\frac{e_01}{2}(s+j)}+\text{ entire},\hfill \end{array}$$
completing the proof.
In particular, if $`A=U(L)`$, then $`\{0\}`$ is a sojourn time if and only if $`L\text{ Lsp}(M,g)`$. If $`L`$ Lsp$`(M,g)`$, $`\zeta (s,L)`$ is regular at 0.
Now let us return to the case where the geodesic flow is completely integrable. Then the dimension of each periodic torus $`𝒯`$ of period $`L`$ equals $`e_o=dim𝒯=n`$. Hence we have:
$$TrU(t)=e_o(t)+\underset{𝒯}{}e_𝒯(t)$$
$`(3.5a)`$
where the sum runs over the periodic tori in $`S^{}M`$ and where
$$e_𝒯(t)=a_{𝒯;\frac{n+1}{2}}(tL+i0)^{\frac{n+1}{2}}+a_{𝒯;\frac{n+1}{2}+1}(tL+i0)^{\frac{n+1}{2}+1}+\mathrm{}.$$
$`(3.5b)`$
More precisely, it takes this form if $`n`$ is even; if $`n`$ is odd, the positive powers of $`(tL+i0)`$ should be multiplied by $`log(tL+i0)`$. In the following Corollarly we use the notation $`𝒯`$ for the time reverse torus $`\sigma 𝒯.`$
(3.6) Corollary
$$\underset{\pm }{}a_{\pm 𝒯,(\frac{n+1}{2})+k}=\text{res}(\sqrt{\mathrm{\Delta }}^{\frac{n+1}{2}+k}U(t)|_{t=L})$$
Proof. In the notation of (possibly negative) fractional derivatives in $`t`$, we have $`\mathrm{\Delta }^\mu U(L)=D_t^\mu U(t)|_{t=L}=`$. The claim follows from the facts that
$$D_t^{\frac{n+1}{2}+k}(tL+i0)^{\frac{n+1}{2}+k}=log(tL+i0)$$
and that the non-commutative residue is the coefficient of $`log(tL+i0)`$.
(3.7) Examples:
(a) If $`dimM=2`$, then the wave trace expansion at a torus $`𝒯`$ has the form
$$e_𝒯(t)=a_{𝒯,\frac{3}{2}}(tL+i0)^{\frac{3}{2}}+a_{𝒯,\frac{1}{2}}(tL+i0)^{\frac{1}{2}}+\mathrm{}.$$
Hence
$$\underset{\pm }{}a_{\pm 𝒯,\frac{3}{2}}=res\sqrt{\mathrm{\Delta }}^{\frac{3}{2}}e^{iL\sqrt{\mathrm{\Delta }}},a_{𝒯,\frac{1}{2}}=res\sqrt{\mathrm{\Delta }}^{\frac{1}{2}}e^{iL\sqrt{\mathrm{\Delta }}},\mathrm{}$$
(b) If $`dimM=3`$, the wave trace expansion at $`𝒯`$ has the form
$$e_𝒯(t)=a_{𝒯,2}(tL+i0)^2+a_{𝒯,1}(tL+i0)^1+a_{𝒯,0}log(tL+i0)+\mathrm{}.$$
Hence
$$\underset{\pm }{}a_{\pm 𝒯,2}=res\mathrm{\Delta }^2e^{iL\sqrt{\mathrm{\Delta }}},a_{𝒯,1}=res\mathrm{\Delta }^1e^{iL\sqrt{\mathrm{\Delta }}},\underset{\pm }{}a_{\pm 𝒯,0}=rese^{iL\sqrt{\mathrm{\Delta }}},\mathrm{}.$$
(3.8) Remark: The wave invariants for a closed geodesic $`\gamma `$ (or periodic torus $`T_L`$) are exactly the same as for its time reversal, hence the same residue formulae also give the individual wave invariants. For this reason it is often not necessary to resolve the ambiguity between a torus and its time reversal.
## 4 Proof of the Main Lemma
We now prove that the wave trace invariants of a metric $`g^{}`$ determine its quantum normal form $`\widehat{H}`$.
### 4.1 Wave invariants for simple surfaces of revolution
We begin by specializing Corollary (3.6) to the case of a surface of revolution $`(S^2,g)`$ in $`^{}.`$ Since the length spectrum $``$ is simple, there is a unique periodic torus $`𝒯_LS_g^{}S^2`$ of each length $`L`$ (up to time reversal). By the existence of global action-angle coordinates, it may be expressed in the form $`𝒯_L=T_{I_L}`$ for a unique point $`I_L=(I_{L1},I_{L2})𝒫`$ (up to reflection). Let $`\omega _L`$ denote the frequency vector at $`I_L`$. Then we have $`L\omega _L=M_L`$ where $`M_L`$ is the vector of winding numbers of the periodic orbits on $`𝒯_L.`$
The other periodic torii of period $`L`$ lie on the rays $`\mathrm{I}\mathrm{R}^+𝒯_L\mathrm{I}\mathrm{R}^+\sigma (𝒯_L)`$. Their action coordinates lie on the rays $`\mathrm{I}\mathrm{R}^+I_L\mathrm{I}\mathrm{R}^+\sigma (I_L)`$ and they have the same frequency vector, $`\omega _L`$, as for $`I_L`$. Since the wave invariants at $`𝒯_L`$ depend only on the microlocalization of the wave group to a conic neighborhood of $`\mathrm{I}\mathrm{R}^+𝒯_L`$, we introduce a microlocal cut off operator $`\widehat{\psi }_L(\widehat{I}_1,\widehat{I}_2)`$ with $`\psi `$ homogeneous of degree 0, equal to 1 in a small conic neighborhood of the ray $`\mathrm{I}\mathrm{R}^+I_L`$ and zero off of a slightly larger conic neighborhood. The singularity of $`TrU(t)`$ at $`t=L`$ is then the same as the singularity of $`Tr\widehat{\psi }_LU(t).`$
To calculate the wave trace as a residue, we also introduce the first order pseudodifferential operator $`\omega _L\widehat{I}:=\omega _{L1}\widehat{I}_1+\omega _{L2}\widehat{I}_2`$. We emphasize that $`\omega _L\widehat{I}`$ is a linear combination with constant coefficients of the action operators. Since $`\omega _II=H`$, the principal symbol $`\omega _LI`$ of $`\omega _L\widehat{I}`$ takes the value 1 at $`I=I_L`$ and therefore $`\omega _L\widehat{I}`$ is elliptic in a conic neighborhood of $`𝒯_L`$. We will use it as the gauging elliptic operator in the residue formula for the wave invariants.
(4.1.1) Proposition Let $`g^{}`$ and let $`L`$. Then we have:
$$\underset{\pm }{}a_{\pm T_L,\frac{1}{2}+k}=Res_{s=0}_\mathrm{\Gamma }\psi _L(I+\mu )e^{iM_L,I}e^{iL\widehat{H}(I+\mu )}(\widehat{H}(I+\mu ))^{\frac{1}{2}+k}(\omega _L(I+\mu ))^s𝑑I$$
where as above $`M_L`$ is the vector of winding numbers of $`𝒯_L`$.
Proof: Since $`\sqrt{\mathrm{\Delta }}=\widehat{H}(\widehat{I}_1,\widehat{I}_2)`$ and since $`\widehat{\psi }_L`$ is a function of the action operators, we have by Corollary (3.6) that
$$\underset{\pm }{}a_{\pm T_L,\frac{1}{2}+k}=Res_{s=0}\underset{NZZ^2}{}\psi _L(N+\mu )e^{iL\widehat{H}(N+\mu )}(\widehat{H}(N+\mu ))^{\frac{1}{2}+k}(\omega _L(N+\mu ))^s.$$
We then apply the Poisson summation formula for Re $`s>>0`$ to replace the sum over $`NZZ^2`$ by
$$Res_{s=0}\underset{MZZ^2}{}J_{L,M,k}(s)$$
where
$$J_{L,M,k}(s):=_\mathrm{\Gamma }\psi _L(I+\mu )e^{iM,I}e^{iL\widehat{H}(I+\mu )}(\widehat{H}(I+\mu ))^{\frac{1}{2}+k}(\omega _L(I+\mu ))^s𝑑I.$$
It follows from Theorem 3.1 and by simplicity of the length spectrum that only the term with $`M=M_L`$ has a pole at $`s=0.`$ This can be seen more directly from the fact that only in this term does the phase $`M,I+LH(I)`$ have a critical point, since $`M=L_IH(I_M)`$ implies that the torus with actions $`I_M`$ is periodic of period $`L`$ (cf. §1.2).
To calculate the residue of the integral $`J_{Lk}(s):=J_{L,M_L,k}(s)`$, we rewrite the integrals in terms of the $`(\rho ,\xi )`$ coordinates introduced in (1.4.3) in a conic neighborhood of the point $`I_L\{H=1\}`$: Recall that we parametrize the tangent line $`T_{I_L}(\{H=1\})`$ by
$$\xi \mathrm{I}\mathrm{R}I_L+\xi v$$
where $`v`$ is a non-zero vector along the line $`M_LI=0`$ and parametrize a conic neighborhood of $`\mathrm{I}\mathrm{R}^+I_L`$ by
$$(\rho ,\xi )\mathrm{I}\mathrm{R}^+\times \mathrm{I}\mathrm{R}\rho (I_L+\xi v).$$
Since $`I_L`$ is fixed, we abbreviate $`h_{I_L}`$ by $`h`$ in the next proposition. We also denote the cutoff in these coordinates by $`\psi _L(\xi ).`$
(4.1.2) Proposition With the above notation, $`a_{T_L,\frac{1}{2}+k}`$ equals the term of order $`\rho ^{1k}`$ in the asymptotic expansion of the integral
$$J_{Lk}(\rho ):=L^{2+sk+\frac{1}{2}}e^{iM,\mu }\rho ^{\frac{1}{2}}_{\mathrm{I}\mathrm{R}}\psi _L(\xi )e^{i\rho \alpha \frac{\xi ^2}{2}}e^{i\rho g_3}e^{i_{k=1}^{\mathrm{}}L^{k+1}\rho ^kh_k(\xi )}[h(\xi )+\underset{\mathrm{}=1}{\overset{\mathrm{}}{}}L^{\mathrm{}+1}\rho ^\mathrm{}1h_{\mathrm{}}(\xi )]^{\frac{1}{2}+k}𝑑\xi .$$
Here, $`\alpha =h^{\prime \prime }(0)`$, and $`g_3`$ is the third order remainder in the Taylor expansion of $`h(\xi )`$ at $`\xi =0.`$
Proof We first change variables $`I+\mu I`$ in the expression for $`J_{L,M_L,k}(s)`$ in the preceding Lemma and then further change variables to $`(\rho :=\omega _MI,\xi )`$. Thus we get
$$J_{Lk}(s)=e^{iM,\mu }_o^{\mathrm{}}_{\mathrm{I}\mathrm{R}}\psi _L(\xi )e^{i\rho M,(I_L+\xi v}e^{iL[\rho h(\xi )+\rho ^1h_1(\xi )+\mathrm{}]}[\rho h(\xi )+\rho ^1h_1(\xi )+\mathrm{}]^{\frac{1}{2}+k}\rho ^{s+1}𝑑\rho 𝑑\xi .$$
Taylor expanding $`h(\xi )=h(0)+h^{}(0)\xi +\frac{1}{2}h^{\prime \prime }(0)\xi ^2+g_3(\xi )`$ and using that $`h(0)=H(I_L)=1`$, that $`h^{}(0)=\omega _{I_L}v=0`$ and that $`M,I_L+\xi v=M,I_L=L`$ by (1.2.1a-b), we get
$$J_{Lk}(s)=e^{iM,\mu }_o^{\mathrm{}}_{\mathrm{I}\mathrm{R}}\psi _L(\xi )e^{iL\rho h^{\prime \prime }(0)\frac{1}{2}\xi ^2}e^{iL\rho g_3(\xi )}e^{iL[\rho ^1h_1(\xi )+\mathrm{}]}[\rho h(\xi )+\rho ^1h_1(\xi )+\mathrm{}]^{\frac{1}{2}+k}\rho ^{s+1}𝑑\rho 𝑑\xi .$$
Now change variables again, $`\rho L^1\rho `$ and pull the factor $`\rho L^1`$ in front of $`h(\xi )`$ in the bracketed expression outside the $`d\xi `$ -integral. We get:
$$J_{Lk}(s)=e^{iM,\mu }L^{2+sk+\frac{1}{2}}_o^{\mathrm{}}\{_{\mathrm{I}\mathrm{R}}\psi _L(\xi )e^{i\rho h^{\prime \prime }(0)\frac{1}{2}\xi ^2}e^{i\rho g_3(\xi )}e^{iL[L\rho ^1h_1(\xi )+\mathrm{}]}$$
$$[h(\xi )+\rho ^2L^2h_1(\xi )+\mathrm{}]^{\frac{1}{2}+k}d\xi \}\rho ^{s+1+k\frac{1}{2}}d\rho .$$
The pole at $`s=0`$ is produced by the terms of order $`\rho ^1`$ in the $`d\xi d\rho `$-integrals, hence by the terms of order $`k1\frac{1}{2}`$ in the asymptotic expansion of the $`d\xi `$-integral.
To determine the terms of order $`\rho ^{k1}`$ in $`J_{Lk}(\rho )`$ we apply the method of stationary phase. Cancelling factors of $`\rho ^{\frac{1}{2}}`$ we get:
(4.1.3) Corollary In the above notation, $`a_{T_L,\frac{1}{2}+k}`$ equals the term of order $`\rho ^{k1}`$ in the asymptotic expansion of
$$\frac{1}{\sqrt{2\pi i\alpha }}L^{2k+\frac{1}{2}}e^{iM,\mu }\underset{m=0}{\overset{\mathrm{}}{}}(2i\rho )^m\alpha ^m_\xi ^{2m}\{e^{i\rho g_3(\xi )}e^{i_{j=1}^{\mathrm{}}L^{j+1}\rho ^jh_j(\xi )}[h(\xi )+\underset{\mathrm{}=1}{\overset{\mathrm{}}{}}L^{\mathrm{}+1}\rho ^\mathrm{}1h_{\mathrm{}}(\xi )]^{\frac{1}{2}+k}\}|_{\xi =0}$$
### 4.2 The meridian torus
We now want to analyse the terms of the expansion coming from the ray of meridian torii with $`I_1=0.`$ The special property of the meridian torii is that they are invariant under the canonical involution $`\sigma (x,\xi )=(x,\xi )`$ of $`T^{}S^20`$. As the following proposition shows, this leads to a useful symmetry property of the quantum normal form. As above, the notation $`h(\xi ),h_1(\xi )`$ etc. refer to the functions $`h_{I_L}`$ (etc.) where $`I_L`$ is the point on $`\{H=1\}`$ corresponding to the meridian torus in $`S_g^{}S^2`$.
(4.2.1) Proposition The complete symbol of $`\widehat{H}`$ is invariant under $`\sigma `$. Hence $`h_j`$ is even for all j.
Proof Let $`C`$ denote the operator of complex conjugation: $`C\psi =\overline{\psi }.`$ Since $`\sqrt{\mathrm{\Delta }}`$ commutes with $`C`$, so does $`\widehat{H}(\widehat{I}).`$ Since $`C`$ is a conjugate-linear involution, this implies that
$$\widehat{H}(\widehat{I})=C^1\widehat{H}(\widehat{I})C=\overline{\widehat{H}}(C^1\widehat{I}C)$$
where the bar denotes complex conjugation.
We claim that $`C^1\widehat{I}_1C=\widehat{I}_1`$ and that $`C^1\widehat{I}_2C=\widehat{I}_1`$. Moreover, that $`I_1\sigma =I_1,I_2\sigma =I_2.`$ The statements regarding $`I_1,\widehat{I}_1`$ are obvious since $`\widehat{I}_1=\frac{}{i\theta }`$ changes sign under complex conjugation and since $`\sigma _{C^1\widehat{I}_1C}=I_1\sigma .`$ This latter also follows from the fact $`I_1(x,\xi )=x_2\xi _1x_1\xi _2`$.
Regarding $`I_2`$ we note that its $`\sigma `$-invariance follows immediately from the explicit formula (1.3.1). The non-obvious claim is that $`\widehat{I}_2`$ is invariant under $`C`$. But from the invariance of the principal symbol we have
$$C^1\widehat{I}_2C=\widehat{I}_2+K_1$$
where $`K_1`$ is of order 0. Since $`\widehat{I}_2`$ is a function of $`(\widehat{I}_1,\sqrt{\mathrm{\Delta }})`$, it is clear that $`C`$ must take joint $`(\widehat{I}_1,\widehat{I}_2)`$ eigenfunctions into joint eigenfunctions. Hence $`C`$ determines an involution on the joint spectrum $`\{(n,k+\frac{1}{2}):|n|2k+1,k0\}`$ which, we recall, is simple. Consequently, the involution (still denoted $`C`$) must take the form
$$C:(n,k+\frac{1}{2})(n,k+\frac{1}{2}+f(n,k))$$
where $`f`$ is a bounded function. Moreover, since $`f(n,k)=\widehat{I}_2\overline{\varphi _{n,k}},\overline{\varphi _{n,k}}`$ and since the $`\varphi _{n,k}`$ are quasi-modes associated to Bohr-Sommerfeld-Maslov torii \[CV.3\], it is clear that $`f(n,k+\frac{1}{2})`$ must be a polyhomogeneous function of order 0 on $`\mathrm{\Gamma }_o`$. Since it is also integral-valued on the semi-lattice of joint spectral points, it must be constant. Additionally it must satisfy the involution condition, and one sees that the constant must be 0.
Returning to $`\widehat{H}(\widehat{I}_1,\widehat{I}_2)`$ we have $`\widehat{H}(\widehat{I}_1,\widehat{I}_2)=C^1\widehat{H}(\widehat{I}_1,\widehat{I}_2)C=\overline{\widehat{H}}(C^1\widehat{I}_1C,C^1\widehat{I}_2C)=\overline{\widehat{H}}(\widehat{I}_1,\overline{I}_2)`$ so that $`\widehat{H}(a,b)=\overline{\widehat{H}}(a,b).`$ We then observe that $`\widehat{H}`$ is a real function (at least modulo terms of order $`\mathrm{}.`$) To see this, we recall that the eigenvalues $`\widehat{H}(n,k+\frac{1}{2})`$ are real and argue by induction on the symbol expansion using that
$$\widehat{H}(n,k+\frac{1}{2})=H(n,k+\frac{1}{2})+H_1(n,k+\frac{1}{2})+\mathrm{}\mathrm{I}\mathrm{R}.$$
First, the principal symbol $`H`$ is real so we may drop it from the expansion without affecting reality. Assuming that $`H=H_1,H_1,\mathrm{},H_k`$ are real, we may drop them all and get that the tail sum is real. Since it is dominated by $`H_{k1}`$, this function must be real for all points $`(n,k+\frac{1}{2})`$ sufficiently far from $`0`$. But by homogeneity, $`H_{k1}`$ is then real on the projection of these lattice points to $`H=1.`$ The projected points form a dense set by Proposition 1.2.4 and hence $`H_{k1}`$ is everywhere real. We conclude that $`\widehat{H}`$ is $`\sigma `$-invariant.
Now consider the ray of meridian torii, or more precisely the ray $`\mathrm{I}\mathrm{R}^+(0,1)`$ in the action cone $`\mathrm{\Gamma }.`$ We note that this ray is invariant under the involution $`\sigma _\mathrm{\Gamma }(a,b):=(a,b)`$ of $`\mathrm{\Gamma }`$, and moreover the meridian torus is invariant under $`\sigma `$ (it ‘rotates’ the torus by angle $`\pi .`$) From the above, the level set $`\{H=1\}`$ is invariant under $`\sigma _\mathrm{\Gamma }`$ and hence the tangent line at $`(0,1)`$ is invariant. Evidently it is horizontal in the $`(I_1,I_2)`$-plane and $`\sigma _\mathrm{\Gamma }`$ restricts to it to the map $`\xi \xi .`$ Since the complete symbol of $`\widehat{H}`$ is $`\sigma _\mathrm{\Gamma }`$-invariant, the $`h_j`$’s must be even.
We now go back to the wave invariants associated to the meridian torus and its iterates. Let us write the amplitude for the kth wave invariant of the iterate of length $`L`$, namely
$$A_{Lk}(\xi ,\rho ):=e^{i_{j=1}^{\mathrm{}}L^{j+1}\rho ^jh_j(\xi )}[h(\xi )+\underset{\mathrm{}=1}{\overset{\mathrm{}}{}}L^{\mathrm{}+1}\rho ^\mathrm{}1h_{\mathrm{}}(\xi )]^{\frac{1}{2}+k},$$
$`(\mathrm{4.2.2})`$
in the form
$$A_{Lk}(\xi ,\rho ):=A_{Lko}(\xi )+\rho ^1A_{Lk1}(\xi )+\mathrm{}.$$
Thus, $`a_{T_L,\frac{1}{2}+k}`$ is the term of order $`\rho ^{k1}`$ in
$$\frac{L^{k1}}{\sqrt{2\pi iL\alpha }}\underset{m=0}{\overset{\mathrm{}}{}}\underset{\mathrm{}=0}{\overset{\mathrm{}}{}}\alpha ^m(2i\rho )^m\mathrm{}_\xi ^{2m}[e^{i\rho g_3}A_{Lk\mathrm{}}].$$
$`(\mathrm{4.2.3})`$
Expanding the derivatives and using that $`g_3(\xi )`$ is even and of order $`0(\xi ^4)`$, we may rewrite (4.2.3) in the form
$$\frac{L^{k1}}{\sqrt{2\pi iL\alpha }}\underset{m=0}{\overset{\mathrm{}}{}}\underset{\mathrm{}=0}{\overset{\mathrm{}}{}}\underset{j=0}{\overset{m}{}}\underset{q\frac{1}{2}m}{}\underset{(i_1,\mathrm{},i_q):|i|=j}{}\{C_{mli}\alpha ^m\rho ^{m\mathrm{}+q}(_\xi ^{2i_1}g_3\mathrm{}_\xi ^{2i_q}g_3)_\xi ^{2(mj)}A_{Lk\mathrm{}}\}.$$
$`(\mathrm{4.2.4})`$
Here, the multi-index $`i`$-sum runs over $`q`$-tuplets with $`i_n2,|i|=i_n=jm`$ with $`q\frac{1}{2}m,2q|i|=jm.`$ In order that $`m\mathrm{}+q=k1`$ it is necessary that
$$m2(k+1),\mathrm{}k+1.$$
These restrictions follow from the fact that it takes at least 4 derivatives of $`g_3`$ to make a non-zero contribution.
### 4.3 A collection of formulae
For future reference we assemble some notation and identities regarding the coefficients $`A_{Lk\mathrm{}}`$ and their relations to the wave invariants.
(4.3.1) Notation
| (1) Define: $`F_n(h_1,\mathrm{},h_n)`$ by $`[h(\xi )+_{\mathrm{}=1}^{\mathrm{}}L^{\mathrm{}+1}\rho ^\mathrm{}1h_{\mathrm{}}(\xi )]^{\frac{1}{2}}=`$ |
| --- |
| $`h(\xi )^{\frac{1}{2}}[1+_{n=2}^{\mathrm{}}L^n\rho ^nF_n(h^1,h_1,\mathrm{},h_n)]`$ |
| (2) Define: $`G_n(L,h_1,\mathrm{},h_n)`$ by $`e^{i_{j=1}^{\mathrm{}}L^{j+1}\rho ^jh_j}=1+_{n=1}^{\mathrm{}}L^n\rho ^nG_n(L,h_1,\mathrm{},h_n).`$ |
| (3) Define $`F_{nmj}`$ by $`F_n(h^1,h_1,\mathrm{},h_n)=_{m=1}^nh^m[_{j=(j_1,\mathrm{},j_m):|j|=n}F_{nmj}h_{j_1}\mathrm{}h_{j_m}]`$ |
| (4) Define $`G_{nmj}`$ by $`G_n(L,h_1,\mathrm{},h_n)=_{m=1}^nL^m[_{j=(j_1,\mathrm{},j_m):|j|=n}G_{nmj}h_{j_1}\mathrm{}h_{j_m}]`$ |
(4.3.2) Identities
| (1) $`A_{L0n}`$ = $`L^nh^{\frac{1}{2}}[F_n+G_n+_{i+j=n}F_iG_j]`$ |
| --- |
| (2) $`A_{Lk}`$ = $`[h(\xi )+_{\mathrm{}=1}^{\mathrm{}}L^{\mathrm{}+1}\rho ^\mathrm{}1h_{\mathrm{}}(\xi )]^kA_{L0}.`$ |
| (3) $`A_{L0n}=L^nh^{\frac{1}{2}}[_{m=1}^n_{j=(j_1,\mathrm{},j_m):|j|=n}\{F_{nmj}h^m+G_{nmj}L^m\}h_{j_1}\mathrm{}h_{j_m}`$ |
| $`+_{a+b=n}_{m_1=1}^a_{m_2=1}^b_{i=(i_1,\mathrm{},i_{m_1}):|i|=a}_{j=(j_1,\mathrm{},j_{m_2}):|j|=b}F_{am_1i}G_{bm_2j}h_{i_1}\mathrm{}h_{i_{m_1}}h_{j_1}\mathrm{}h_{j_{m_2}}].`$ |
| (4) $`A_{Lk\mathrm{}}=hA_{L(k1)\mathrm{}}+_{ij:i+j+1=\mathrm{},j1}L^{j+1}h_jA_{Lki}.`$ |
The coefficients $`F_{nmj},G_{nmj}`$ are universal and hence up to the prefactor of $`h^{\frac{1}{2}}`$, $`A_{Lkj}`$ is a (non-homogeneous) polynomial of degree $`j`$ in $`L`$, in $`h^1,`$ and in the $`h_j`$’s with universal coefficients. In particular, the first few $`A_{Lkj}`$’s are given by:
$$A_{Lko}=h(\xi )^{\frac{1}{2}+k},A_{Lk1}=h(\xi )^{\frac{1}{2}+k}L^2h_1,A_{Lk2}=h(\xi )^{\frac{1}{2}+k}[L^3h_2+L^4h_1^2+(k\frac{1}{2})L^2\frac{h_1}{h}].$$
It follows that the principal wave invariant is given by
$$a_{T_L,\frac{3}{2}}:=c_L=\frac{1}{\sqrt{2\pi i\alpha L}}e^{iM_L,\mu }$$
and that the higher wave invariants $`a_{T_L,\frac{1}{2}+k}`$ are given by $`c_L`$ times polynomials in $`L`$ and in the derivatives of $`h,h^1,h_1,\mathrm{}`$ at $`\xi =0`$. For instance, the subprincipal wave invariant in dimension 2 is given given in terms of universal coefficients $`C_{ijk}^{}`$ by:
$$a_{T_L,\frac{1}{2}}=c_L[C_{004}^{}_\xi ^4g_3A_{L00}+C_{010}^{}A_{L01}+C_{002}^{}_\xi ^2A_{L00}]$$
which is easily seen to equal
$$c_L[C_{004}_\xi ^4h(\xi )|_{\xi =0}+C_{010}L^2h_1(0)+C_{002}_\xi ^2h(\xi )^{\frac{1}{2}}|_{\xi =0}].$$
Above the indices in the coefficients $`C_{kmj_1j_2\mathrm{}}`$ have the following meaning: $`k`$ corresponds of course to the $`k`$ index; $`m`$ is the power of $`L`$ and $`j_n`$ are the jet orders of $`h_n`$, with with exception of $`j_o`$ which is the jet order of $`h`$.
(4.3.3) Proposition We have
$$a_{L,\frac{1}{2}+k}=c_LP_k(L,h^{(2)}(0),\mathrm{},h^{(2k+4)}(0),h_1(0),\mathrm{},h_1^{(2k)}(0),\mathrm{},h_k(0),h_k^2(0),h_{k1}(0))$$
where $`P_k`$ is a polynomial with the following properties:
(i) It involves only the first $`2k+4`$ Taylor coefficients of $`h`$ at $`0`$, the first $`2k`$ of $`h_1`$,…, the first $`2k+22n`$ of $`h_n`$ …, the first $`2`$ of $`h_k`$ and the 0th of $`h_{k1}.`$
(ii) It is of degree 1 in the variables $`h_{k1}(0)`$, $`h_k^{(2)}(0),\mathrm{},h^{(2k+4)}(0)`$, and each occurs in precisely one term.
(iii) The $`L`$-order of the monomials containing these terms is respectively $`L^{k+2},L^{k+1},\mathrm{},L^0.`$
Proof: To prove these claims we combine the formula for $`a_{T_L,k+\frac{1}{2}}`$ in (4.2.4) with the formulae for the $`A_{Lk\mathrm{}}`$ given in (4.3.2). Since we may replace the factors of $`g_3`$ by $`h`$ in (4.2.4) and since (4.3.2) expresses $`A_{Lk\mathrm{}}`$ as $`h^{\frac{1}{2}}`$ times a polynomial in $`L`$ and in the $`h_j`$’s, it is clear that $`a_{T_L,k+\frac{1}{2}}`$ is given by a polynomial in the data stated above. It remains to prove that the polynomial has the properties claimed in (i)–(iii).
We prove these claims by proving the stronger statement that
$$a_{L,\frac{1}{2}+k}=c_LQ_k(L,h(\xi )^{\frac{1}{2}},h(\xi ),h^{(2)}(\xi )),\mathrm{},h^{(2k+4)}(\xi )),h_1(\xi )),\mathrm{},h_1^{(2k)}(\xi )),\mathrm{},h_k(\xi )),h_k^2(\xi )),h_{k1}(\xi )))|_{\xi =0}$$
where $`Q_k`$ has the properties (i)-(iii) for variable $`\xi .`$
The proof is by induction on $`k`$. The properties are visibly true for the principal and subprincipal wave invariants. Assume then that they are correct for $`kN`$ and consider how things change as $`NN+1`$. First, the amplitude $`A_{L(N+1)}`$ is given by $`[h+L^2\rho ^2h_1+\mathrm{}]A_{LN}`$. Second, we are looking at the term of order $`\rho ^{N2}`$ rather than $`\rho ^{N1}`$ in the asymptotic series (4.2.4).
In going one further order into the asymptotic series, only two new things happen:
$``$ The term $`h_{N2}`$ appears for the first time, arising from the linear term in the Taylor expansion of the exponential in (4.2.2). The linear term in the binomial expansion of the power in (4.2.2) does not contribute at this stage because its $`\rho ^1`$ factor has one higher power.
$``$ Two additional derivatives are allowed to fall on the previous $`h_j`$’s. A priori, there could be between two and six additional derivatives in a method of stationary phase expansion. However, the cases of three to six derivatives do not contribute any new data. Indeed, the cases of three-four derivatives occur when $`m`$ goes up by one. One has to remove the extra factor of $`\rho ^1`$ by applying at least one derivative to the $`e^{i\rho g_3}`$ factor. But in fact all four have to be applied to $`g_3`$ to make a non-zero contribution, and hence no derivatives are left to apply to the $`h_{\mathrm{}}`$’s. In the case of five or six derivatives, where $`m`$ goes up to two, one needs to remove two extra factors of $`\rho ^1`$ by applying derivatives to the $`e^{i\rho g_3}`$ factor, and there is no non-zero way to do this.
Claim (i) follows immediately from these observations.
To prove (ii) we note that by the induction hypothesis, $`h_{N1}(\xi ),\mathrm{},h^{(2N+4)}(\xi )`$ occur linearly in $`P_N(\xi ).`$ Since it requires both of the two new derivatives to fall on these factors to produce $`h_{N1}^{(2)}`$ etc. in $`P_{N+1}(\xi )`$, these factors will also occur to order 1. As for $`h_{N2}`$, we observed above that it comes only from the linear term in the exponential in (4.2.2) and hence it too appears to order 1.
The proof of (iii) follows the proof of (ii). Since the terms in the exponent of the exponential factor in (4.2.2) have the form $`L^{j+1}\rho ^j`$, the new $`h_{N2}(\xi )`$ term has the coefficient $`L^{N+3}`$. Similarly, the other terms under discussion, e.g. $`h_{N2+r}^{(2r)}(\xi )`$, originated as $`h_{N2+r}`$ at the $`N+1r`$th stage with the factor of $`L^{N+3r}`$. We would like to show that they remain with just this factor of $`L`$ as we move inductively up to the $`N+1`$st stage. To this end, we note that in order to produce the term $`h_{N2r}^{(2r)}`$ at the (N+1)st stage, $`2r`$ derivatives must fall on the original $`h_{N2+r}`$. But there are exactly $`r`$ stages intermediate between the $`N+1r`$th stage and the $`N+1`$st stage and at each stage at most two new derivatives can fall on a factor. Hence, each of the two new derivatives at each stage must fall on the factor of concern.
Let us also consider what can be multiplied against the original $`h_{N2+r}`$ in the course of producing $`h_{N2+r}^{(2r)}(\xi )`$. We observe that each new pair of derivatives is accompanied by a factor of $`\rho ^1.`$ Since $`h_{N2+r}(\xi )`$, at the $`N+1r`$th stage, comes with the factor $`\rho ^{N2+r}`$, the $`2r`$ further derivatives will bring its $`\rho ^1`$-order up to $`\rho ^{N2}.`$ Hence no other factors of $`\rho ^1`$ could have fallen on this factor. Therefore in the repeated multiplications by $`[h+L^2\rho ^2h_1+\mathrm{}]`$, only the repeated choice of the term $`h`$, of order 0 in $`\rho ^1`$, can have given rise to the term $`h_{N2+r}^{(2r)}`$ . It follows that it retains its original $`L`$-order, namely $`L^{N+3r}.`$
### 4.4 Completion of Proof of Main Lemma
Our purpose is now to show that by using the joint $`\rho ,L\mathrm{}`$ asymptotics one can recover the complete Taylor expansions of all the $`h_j`$’s from the wave invariants of the meridian torus and its iterates. That is, to complete the
Proof of the Main Lemma: We will prove by induction that from the wave invariants $`a_{T_{pL},\frac{1}{2}+k}`$ with $`kN`$ and all $`p𝐍`$, we can determine the $`2k+4`$ -jet of $`h`$, the $`2k`$-jet of $`h_1`$,…, the $`2k+22n`$-jet of $`h_n`$ at $`\xi =0.`$ (for $`nk+1.`$)
We note that, unlike in the non-degenerate case, the principal wave invariant determines the 2-jet of $`h`$, and gives no new information under iteration. We therefore begin the induction with the subprincipal wave invariant $`a_{T,\frac{1}{2}}`$.
From the explicit formula above for the subprincipal wave invariant $`a_{L,\frac{1}{2}}`$ it is evident that the 2-jet term in $`h`$ is old information, while the other two terms differ in the power of $`L`$. Hence they decouple under iteration $`LpL`$ and we can determine the 4-jet of $`h`$ and the 0-jet of $`h_1`$ from the first two wave invariants. Hence the induction hypothesis is correct at the first stage.
Assume the induction hypothesis for the $`k1`$st stage. Then the only new information at the kth stage is that contained in the terms $`h_{k1}(0)`$, $`h_k^{(2)}(0),\mathrm{},h^{(2k+4)}(0)`$. By the proposition above, they occur linearly in the monomials containing them and the monomials have also the factor of $`L`$ to the powers $`L^{k+1},L^k,\mathrm{}.`$ Hence the terms containing the new data decouple as $`p\mathrm{}`$, and the new data can be determined as stated. This completes the inductive argument.
It follows that we can determine the full Taylor expansions of the $`h_j`$’s at $`\xi =0.`$ Since they are real analytic they are completely determined. Then from the homogeneity of $`H_j`$, we can determine $`H_j`$ from $`h_j`$ and hence the entire function $`\widehat{H}`$ is determined.
Remark In the above argument, we are able to drop the terms involving only previously known data because of the universal nature of the polynomials $`P_k`$. This universality depends on the fact that we are only comparing wave data for quantum torus integrable Laplacians.
## 5 Proofs of Corollaries 1 and 2
We now show that the symplectic equivalence class of a metric in $`^{}`$ is spectrally determined among metrics in this class.
(5.1) Proof of Corollary 1 By the Main Lemma, the principal symbol $`H_g(I_1,I_2)`$ as a function of the action variables is spectrally determined for a metric $`g`$ in $`^{}.`$ By Propositions (1.5.2-3), $`H_g`$ determines the geodesic flow up to symplectic equivalence among metrics in $`^{}`$.
Remark The Corollary could also be proved by noting that the first return time $`\tau _E`$ is spectrally determined. But geodesic flows of in the class $`^{}`$ are symplectically equivalent if and only if they have the same first return times $`\tau _E`$. (See also \[C.K\] for an equivalent statement.) Also, it should be noted that the first return time could be determined from the wave invariants at the equator; hence Corollary 1 would also follow from Guillemin’s inverse result for non-degenerate elliptic closed geodesics.
(5.2) Proof of Corollary 2: Since $`\tau _E(I)`$ is spectrally determined, so is the function $`s(i(I))`$ of Proposition 1.3.6. It is given by
$$s(i(I))=_{i(I)}^{\pi i(I)}f(cosu)sin(u)(sin^2usin^2(i(I)))_+^{\frac{1}{2}}𝑑u$$
or, putting $`x=cosu`$ for $`u[0,\frac{\pi }{2}]`$ and $`x=cosu`$ for $`u[\frac{\pi }{2},\pi ]`$, by
$$s(i(I))=_0^{cos(i(I))}[f(x)+f(x)](cos(i(I))^2x^2)_+^{\frac{1}{2}}𝑑x.$$
It therefore suffices to show that
$$Tf(u)=_0^u[f(x)+f(x)](u^2x^2)_+^{\frac{1}{2}}𝑑x$$
determines $`[f(x)+f(x)]`$. But $`[f(x)+f(x)]`$ is smooth and even so may be written as $`g(x^2)`$ for a smooth $`g`$; changing variables $`y=x^2`$, and $`v=u^2`$ we get
$$Tf(v)=_o^1g(y)y^{\frac{1}{2}}(yv)_+^{\frac{1}{2}}𝑑y.$$
Thus $`T`$ is a standard Abel transform, which is well-known to be invertible. It follows that $`g(y)y^{\frac{1}{2}}`$ is spectrally determined; hence so is the even part of $`f`$.
## 6 Proof of Final Lemma
To complete the proof of the Theorem, we need to show that $`\widehat{H}`$ determines a metric in $`.`$ It is plausible that this can be done for the following reason: The spectrum of $`\mathrm{\Delta }`$ is the set of values $`\{\widehat{H}(n,k+\frac{1}{2})^2\}`$ of $`\widehat{H}^2`$ on the integer points of the action cone $`\mathrm{\Gamma }.`$ On the other hand since $`\widehat{I}_1=\frac{}{\theta }`$, the set of values $`\widehat{H}(n,k+\frac{1}{2})`$ for fixed $`n`$, is just the spectrum $`\{E_{nk}\}`$ of the the singular Sturm-Liouville operator
$$L_n=(\frac{d}{dr})^2+q_n(r),q_n(r):=q(r)+\frac{n^2}{a(r)^2},q(r)=\frac{2a(r)a^{\prime \prime }(r)(a^{}(r))^2}{a(r)^2}$$
obtained by separating variables in $`\mathrm{\Delta }`$, fixing the eigenvalue of $`\frac{}{i\theta }`$ equal to $`n`$, and putting the radial part in normal form. Hence, from the coefficients of $`\widehat{H}`$ we can determine $`Spec(L_n)`$ for each $`n`$. That is, from $`Spec(\mathrm{\Delta })`$ we have determined the joint spectrum of $`(\mathrm{\Delta },\frac{}{\theta }).`$ It would thus remain to show that the metric can be determined from this joint spectrum, a much more elementary inverse result which has been stated several times in the literature (\[Kac\],\[B\] \[Gur\]). Since these prior discussions seem to us somewhat sketchy and incomplete, we give a self-contained proof below which was found before we were aware of these references.
Proof of Final Lemma
The proof is basically to write down explicit expressions for $`H`$ and $`H_1`$ in terms of the metric (i.e. in terms of $`a(r)`$) and then to invert the expressions to determine $`a(r)`$. The first step is therefore to construct an initial part of the quantum normal form explicitly from the metric. Up to now, we only know that a polyhomogeneous normal form exists.
To be sure, the principal term $`H`$ has already been implicitly expressed in terms of the metric: Knowledge of $`H`$ is equivalent to knowledge of the level set $`\{H=1\}`$ and hence to knowledge of the function $`F(I_1)`$ in (1.3.1). Unfortunately, we have seen in Corollary (5.2) that $`H`$ only determines the ‘even part’ of the metric. Hence we need to determine at least one of the subprincipal terms $`H_j`$. It turns out that only $`H_1`$ is needed in addition to $`H`$ to determine $`g`$. Since the calculation of $`H_1`$ requires a new calculation of $`H`$, we start calculating both from scratch.
To determine $`H`$ and $`H_1`$ in terms of $`g`$, we are going to study the spectral asymptotics of $`\sqrt{\mathrm{\Delta }}=\widehat{H}(\widehat{I}_1,\widehat{I_2})`$ along ‘rays of representations’ of the quantum torus action, i.e. along multiples of a given lattice point $`(n_o,k_o).`$ Such rays are the quantum analogue of rays $`\mathrm{I}\mathrm{R}^+T_IT^{}S^20`$ thru invariant tori and are basic to homogeneous quantization theory \[G.S.3\]. The basic idea is that the lattice points $`(n_o,k_o+\frac{1}{2})`$ parametrize tori $`T_{n_o,k_o}`$ satisfying the Bohr-Sommerfeld quantization condition. To each such quantizable torus one can construct a joint eigenfunction $`\varphi _{n_o,k_o}`$ of $`(\widehat{I}_1,\widehat{I}_2)`$ by the WKB method. The $`\{\varphi _{n,k}\}`$ are eigenfunctions of $`\mathrm{\Delta }`$ with complete asymptotic expansions along rays. By studying the eigenvalue problem as $`|(n,k)|\mathrm{}`$ we can determine the $`H_j`$’s.
The WKB method we employ is closely related to the classical WKB method for constructing quasi-modes (cf. \[CV.3\] and the Appendix), except that we have an internal rather than an external Planck constant. Let us recall the relevant terminology and notation from the latter case. For each torus $`T_I`$, we denote by $`𝒪^\mu (T_I,A)`$ the space of oscillatory integrals (semi-classical Lagrangean distributions) associated to $`T_I`$, with semi- classical parameters $`\{k_m\}`$ in a set $`A\mathrm{I}\mathrm{R}^+`$ to be specified by the quantization condition. Such oscillatory integrals have the form
$$u(x,k)=k^{\frac{n}{2}+\mu }\underset{\mathrm{}L}{}_𝒱_{\mathrm{}}e^{ik\psi _{\mathrm{}}(x,\xi )}\alpha (x,\xi ,k)𝑑\xi $$
$`(6.1)`$
where the projection of $`T_I`$ is covered by open sets $`𝒰_{\mathrm{}}`$, where $`T^{}(𝒰_{\mathrm{}})𝒰_{\mathrm{}}\times 𝒱_{\mathrm{}}`$ and where the phase $`\varphi _{\mathrm{}}(x,\xi ),`$ with $`(x,\xi )𝒰_{\mathrm{}}\times 𝒱_{\mathrm{}}`$, parametrizes a part of $`T_I`$. The amplitude is a classical symbol in $`k`$ of order $`\mu `$. For further details, we refer to \[C.V.3, §8\].
As discussed above, the torii $`T_I`$ project $`21`$ to the annuli $`r_+(I)rr_{}(I)`$ in $`S^2`$ and have fold singularities along the extremal parallels. Away from the parallels, an associated quasi-mode is given by a sum of two simple WKB functions $`\alpha _\pm (r,\theta )(x)e^{iks\psi _\pm (r,\theta )}.`$ Since the actual $`\mathrm{\Delta }`$\- eigenfunctions $`\varphi _{n_o,k_o}`$ are quasi-modes attached to the Bohr-Sommerfeld torii $`T_{n_ok_o}`$, they have such a form modulo $`|(n_o,k_o)|^{\mathrm{}}.`$ And, since $`\varphi _{n_o,k_o}`$ is an exact $`\frac{}{\theta }`$-eigenfunction, its phases must take the form $`\varphi _{\mathrm{}}(r,\theta )=n_o\theta +S_{n_ok_o}(r)`$ with $`\theta `$-independent amplitudes in polar coordinates. It follows that $`\varphi _{n_o,k_o}(r,\theta )`$ has the form $`e^{in_o\theta }f_{(n_o,k_o)}(r)`$, where $`f_{(n_o,k_o)}(r)`$ is an oscillatory integral in the $`r`$-variable. It is of course associated to the pushed-forward Lagrangean $`\mathrm{\Lambda }_{n_ok_o}:=p_{}T_{n_ok_o}`$ where $`p:T^{}S^2T^{}[0,L]`$ is the projection induced from the map $`(r,\theta )r.`$ In the case of the meridian torii $`T_{0,k_o}`$, the pushforward is just a horizontal line $`p_r=C`$ in $`T^{}[0,L]`$. In the other cases, the $`\mathrm{\Lambda }_{n_ok_o}`$ is a closed curve projecting to an interval $`[r_+(n_ok_o),r_{}(n_ok_o)]`$ with fold singularities at the turning points (endpoints). The curve is given by an equation of the form $`H_{n_o}(r,p_r)=E`$ where $`H_{n_o}(r)=p_r^2\frac{n_o^2}{a(r)^2}`$ is the radial Hamiltonian and where $`E`$ was the level of the torus.
The radial part $`f_{n_o,k_o}(r)`$ of $`\varphi _{n_o,k_o}`$, is an eigenfunction of the radial operator $`D_{n_o}=(\frac{d}{dr})^2\frac{a^{}}{a}\frac{d}{dr}+\frac{n_o^2}{a(r)^2}`$ arising from separating variables in $`\mathrm{\Delta }`$. Before proceeding, we simplify by conjugating $`D_{n_o}`$ to the 1/2-density radial Laplacian
$$\widehat{D}_{n_o}:=a(r)^{\frac{1}{2}}D_{n_o}a(r)^{\frac{1}{2}}=D_r^2+\frac{n_o^2}{a(r)^2}+W$$
where $`W=a(r)^{\frac{1}{2}}[(\frac{d}{dr})^2\frac{a^{}}{a}\frac{d}{dr}]a(r)^{\frac{1}{2}}.`$ (Note that the volume form $`dv_g=a(r)drd\theta `$ of $`(S^2,g)`$ projects to $`a(r)dr.`$) Thus, we view the radial eigenfunction as having the form $`g_{n_o,k_o}(r)\sqrt{adr}`$ and apply the WKB method to the eigenvalue problem
$$\widehat{D}_{n_o}g_{n_o,k_o}(r)(H(n_o,k_o+\frac{1}{2})+H_1(n_o,k_o+\frac{1}{2})+\mathrm{})^2g_{n_o,k_o}(r).$$
$`(6.2)`$
Although the coefficients become singular at $`r=0,L`$, the standard WKB theory applies in the interior because it applies to quasi-modes of $`\mathrm{\Delta }`$ on $`S^2`$. Away from the turning points, the radial part of the 1/2-density eigenfunction therefore has the form
$$g_{n_o,k_o}=\underset{\pm }{}[e^{\pm iS_{n_ok_o}(r)}\underset{m=0}{\overset{\mathrm{}}{}}\alpha _{n_o,k_o;m}(r)]$$
$`(6.3)`$
where the phase $`S_{n_o,k_o}`$ is homogeneous of degree 1 and where the amplitude $`\alpha _{n_o,k_o;m}`$ is homogeneous of degree $`m`$ in $`(n_o,k_o+\frac{1}{2}).`$ These homogeneities replace the powers of $`k`$ in the non-homogeneous theory described above.
The Bohr-Sommerfeld quantization condition on $`T_I`$ is that:
$$\frac{1}{2\pi }I=(n_o,k_o)+\frac{1}{4}\mu _o,(n,k)\mathrm{\Gamma }ZZ^2$$
where $`\mu _0=(0,2)`$ is the Maslov index (cf. \[CV.3, §4\]). It is satisfied by $`T_{n_o,k_o}`$ and hence by the radial Lagrangean $`\mathrm{\Lambda }_{n_o,k_o}`$. The local WKB ansatz (6.3) therefore extends to a quasi-mode of infinite order associated to the global $`\mathrm{\Lambda }_{n_o,k_o}`$. Our purpose now is to write down and solve the first two transport equations, which are needed to determine $`H`$ and $`H_1`$. For background on the relevant aspects of the WKB method, we refer to the Appendix.
The transport equations are obtained by separating out terms of like order in the asymptotic eigenvalue problem
$$\widehat{D}_{n_o}\underset{\pm }{}[e^{\pm S_{n_ok_o}(r)}\underset{m=0}{\overset{\mathrm{}}{}}\alpha _{n_o,k_o;m}(r)](H(n_o,k_o+\frac{1}{2})+H_1(n_o,k_o+\frac{1}{2})+\mathrm{})^2\underset{\pm }{}[e^{\pm S_{n_ok_o}(r)}\underset{m=0}{\overset{\mathrm{}}{}}\alpha _{n_o,k_o;m}(r)].$$
The leading term, of order 2, is the eikonal equation
$$|S_{n_ok_o}^{}(r)|^2+\frac{n_o^2}{a(r)^2}=H(n_o,k_o+\frac{1}{2})^2$$
$`(6.4a)`$
whose solution is the first order phase function
$$S_{n_ok_o}(r):=\sqrt{H(n_o,k_o)^2\frac{n_o^2}{a(r)^2}}𝑑r.$$
$`(6.4b)`$
The Bohr-Sommerfeld quantization condition on $`\mathrm{\Lambda }_{(n_o,k_o)}`$ thus reads:
$$I_{n_o}(E):=\text{Area}\{H_{n_o}E\}=2\pi (k_o+\frac{1}{2})\text{with}E=H(n_o,k_o+\frac{1}{2}).$$
The first transport equation, $`2\alpha _{n_o,k_o;0}^{}S^{}+\alpha S^{}=0`$, is solved by
$$\alpha _{(n_o,k_o;0}(r)=[E_{n_o,k_o}\frac{1}{a(r)^2}]^{\frac{1}{4}}$$
$`(6.5)`$
away from the turning points. Here, $`E_{n_o,k_o}=\frac{H(n_o,k_o+\frac{1}{2})^2}{n_o^2}`$. The solution is of course determined only up to a constant, and we have normalized it so that $`\alpha _{(n_o,k_o;0}`$ is homogeneous of order 1.
In the usual way (see the Appendix), we will interpret $`\alpha _{n_o,k_o;0}`$ as the coefficient of the $`\mathrm{\Xi }_{H_{n_o}}`$\- invariant 1/2-density
$$\nu _0=[E_{n_o,k_o}\frac{1}{a(r)^2}]^{\frac{1}{4}}\sqrt{dr}$$
on $`\mathrm{\Lambda }_{n_ok_o}=\{H_{n_o}(r,p_r)=H(n_o,k_o+\frac{1}{2})\},`$ where $`\mathrm{\Xi }_{H_{n_o}}`$ is the Hamilton vector field of $`H_{n_o}`$. We then re-write the higher coefficients $`\alpha _{n_o,k_o;m}(r)`$ in the form $`\alpha _{n_o,k_o;m}(r)\nu _0.`$
The second transport equation (of order zero) then has the form:
$$\mathrm{\Xi }_{H_{n_o}}\alpha _{n_o,k_o,1}=\alpha _{n_o,k_o,0}^1(\frac{d}{dr})^2\alpha _{n_o,k_o,0}+W+2H(n_o,k_o+\frac{1}{2})H_1(n_o,k_o+\frac{1}{2}),$$
$`(6.6)`$
where $`r`$ denotes the local coordinate on $`\mathrm{\Lambda }_{n_ok_o}`$ obtained by pulling back the base coordinate under the projection. The integral of the left hand side over the closed curve $`\mathrm{\Lambda }_{n_ok_o}=\{H_{n_o}=H(n_o,k_o+\frac{1}{2})\}`$ with respect to the $`\mathrm{\Xi }_{H_{n_o}}`$-invariant density
$$\alpha _{n_o,k_o,0}^2dr=\frac{1}{\sqrt{E_{n_o,k_o}\frac{1}{a(r)^2}}}dr$$
must equal zero. This gives a formula for the first correction to the Bohr-Sommerfeld eigenvalue $`H(k_o,n_o)`$:
$$H_1(n_o,k_o+\frac{1}{2})=\frac{1}{2T(n_o,k_o)H(n_o,k_o+\frac{1}{2})}_{r_{}(n_o,k_o)}^{r_+(n_o,k_o)}[\alpha _{n_o,k_o,0}(\frac{d}{dr})^2\alpha _{n_o,k_o,0}\frac{1}{2}W\alpha _{n_o,k_o,0}^2]𝑑r.$$
$`(6.7)`$
Here, $`T(n_o,k_o)`$ denotes the period of the $`\mathrm{\Xi }_{H_{n_o}}`$ flow on $`\mathrm{\Lambda }_{n_ok_o}`$ and $`r_\pm (n_o,k_o)`$ are the turning points. Plugging in (6.5) we get, formally, the expression
$$H_1(n_o,k_o+\frac{1}{2})=\frac{1}{2T(n_o,k_o)H(n_o,k_o+\frac{1}{2})}_{r_{}(n_o,k_o)}^{r_+(n_o,k_o)}(E_{n_o,k_o}\frac{1}{a(r)^2})^{1/4}(\frac{d}{dr})^2(E_{n_o,k_o}\frac{1}{a(r)^2})^{1/4}𝑑r$$
$`(6.8)`$
$$\frac{1}{2T(n_o,k_o)H(n_o,k_o+\frac{1}{2})}_{r_{}(n_o,k_o)}^{r_+(n_o,k_o)}\frac{W}{\sqrt{E_{n_o,k_o}\frac{1}{a(r)^2}}}𝑑r.$$
Actually, the integral in (6.8) diverges at the turning points, and the correct formula is the regularization obtained (for instance) by the method of the Maslov canonical operator. For the sake of completeness, we provide an exposition of this method in the appendix. Roughly speaking, it regularizes (6.8) by formally integrating the $`\frac{d}{dr}`$ derivatives by parts and by moving them outside the integral (in the appropriate way) as derivatives in the energy level $`E`$. A crucial consequence of the regularization is that only first derivatives of $`a`$ appear in the formulae for $`H_1(n_o,k_o+\frac{1}{2}).`$
Thus the first term of (6.8) is regularized by
$$\frac{C_1}{T(n_o,k_o)H(n_o,k_o+\frac{1}{2})}_E^2_{r_{}(n_o,k_o)}^{r_+(n_o,k_o)}\frac{a^{}(r)^2}{a(r)^6}(E\frac{1}{a(r)^2})^{\frac{1}{2}}𝑑r|_{E=E_{(n_o,k_o)}}.$$
$`(\mathrm{6.8.1}reg)`$
Here and below $`C_i`$ denote (non-zero) constants which we will not need to determine. For the second term of (6.8), we note that
$$_{r_{}(n_o,k_o)}^{r_+(n_o,k_o)}W\alpha _{n_o,k_o;0}^2=_{r_{}(n_o,k_o)}^{r_+(n_o,k_o)}\frac{d}{dr}(\alpha _{n_o,k_o;0}^2a^{\frac{1}{2}})\frac{d}{dr}a^{\frac{1}{2}}𝑑r+_{r_{}(n_o,k_o)}^{r_+(n_o,k_o)}\alpha _{n_o,k_o;0}^2a^{\frac{1}{2}}\frac{a^{}}{a}\frac{d}{dr}a^{\frac{1}{2}}𝑑r.$$
After some simplification, this regularizes to:
$$\frac{C_2}{T(n_o,k_o)H(n_o,k_o+\frac{1}{2})}_E_{r_{}(n_o,k_o)}^{r_+(n_o,k_o)}\frac{a^{}(r)^2}{a(r)^4}(E\frac{1}{a(r)^2})^{\frac{1}{2}}𝑑r|_{E=E_{(n_o,k_o)}}$$
$`(\mathrm{6.8.2}reg)`$
$$+\frac{C_3}{T(n_o,k_o)H(n_o,k_o+\frac{1}{2})}_{r_{}(n_o,k_o)}^{r_+(n_o,k_o)}\frac{a^{}(r)^2}{a(r)^2}(E\frac{1}{a(r)^2})^{\frac{1}{2}}𝑑r|_{E=E_{(n_o,k_o)}}.$$
Now let us return to the inverse problem. We begin from the fact that $`\widehat{H}(n_o,\widehat{I}_2)`$ is a known function of the variable $`I:=\widehat{I}_2`$ for each $`n_o`$. Its principal symbol $`H_{n_o}(I):=H_1(n_o,I)`$ is then a known function and hence the inverse function
$$I_{n_o}(E)=_{r_{}(E)}^{r_+(E)}\sqrt{E\frac{1}{a(r)^2}}𝑑r$$
satisfying $`H_{n_o}(I_{n_o}(E))=E`$ is a known of $`E`$. This of course presupposes that $`_IH_{n_o}(I)0`$, which follows from the non-degeneracy assumption (1.1.6). We may write the integral in the form
$$_{\mathrm{I}\mathrm{R}}(Ex)_+^{\frac{1}{2}}𝑑\mu (x)$$
where $`\mu `$ is the distribution function of $`\frac{1}{a^2}`$, i.e.
$$\mu (x):=|\{r:\frac{1}{a(r)^2}x\}|$$
with $`||`$ the Lebesgue measure. The above integral is an Abel transform and as mentioned above it is invertible. Hence
$$d\mu (x)=\underset{r:\frac{1}{a(r)^2}=x}{}|\frac{d}{dr}\frac{1}{a(r)^2}|^1dx$$
is a spectral invariant. Some simplification leads to the conclusion that the function
$$J(x):=\underset{r:a(r)=x}{}\frac{1}{|a^{}(r)|}$$
is known from the spectrum. By the simplicity assumption on $`a`$, there are just two solutions of $`a(r)=x`$; the smaller will be written $`r_{}(x)`$ and the larger, $`r_+(x).`$ Thus, the function
$$J(x)=\frac{1}{|a^{}(r_{}(x))|}+\frac{1}{|a^{}(r_+(x))|}$$
$`(6.9)`$
is a spectral invariant.
Now let us turn to the $`H_1`$ expression. Since $`H(I_1,I_2)`$ is a spectral invariant, the factors $`H(n_o,k_o+\frac{1}{2})`$ and $`T(n_o,k_o)`$ are spectral invariants. Hence we may remove them from the expression for $`H_1`$ and still get a spectral invariant. For various universal constants $`C_1,C_2,C_3`$ it takes the form
$$[C_1_E^2\frac{(a^{})^2}{a^6}(E\frac{1}{a^2})_+^{\frac{1}{2}}dr+C_2_E\frac{(a^{})^2}{a^4}(E\frac{1}{a^2})_+^{\frac{1}{2}}dr$$
$`(6.10)`$
$$+C_3\frac{(a^{})^2}{a^2}(E\frac{1}{a^2})_+^{\frac{1}{2}}dr]|_{E=\frac{H(I_1,I_2)}{I_1^2}}.$$
By a change of variables, we may rewrite (6.10) in the form
$$[C_1_E^2K(x)x^{\frac{3}{2}}(Ex)_+^{\frac{1}{2}}dx+C_2_EK(x)x^{\frac{1}{2}}(Ex)_+^{\frac{1}{2}}dx$$
$`(6.11)`$
$$+C_3K(x)x^{\frac{1}{2}}(Ex)_+^{\frac{1}{2}}dx]|_{E=\frac{H(I_1,I_2)}{I_1^2}}$$
where
$$K(x)=|a^{}(r_{}(x))|+|a^{}(r_+(x))|.$$
$`(6.12)`$
All values of $`E`$ which occur as ratios $`\frac{H(I_1,I_2)}{I_1^2}`$ give spectral invariants, so (1.10) (as a function of the variable $`E`$) is a spectral invariant.
We now claim that $`K`$ itself is a spectral invariant. To determine it from (6.11) we rewrite (6.11) in terms of the fractional integral operators (cf. \[G.Sh, Ch.1 §5.5\])
$$I_\alpha f(E)=f\frac{x_+^{\alpha 1}}{\mathrm{\Gamma }(\alpha )}(E)=\frac{1}{\mathrm{\Gamma }(\alpha )}_0^Ef(y)(Ey)^{\alpha 1}𝑑y$$
on the half-line $`[0,\mathrm{}].`$ These operators satisfy
$$I_\alpha I_\beta =I_{\alpha +\beta },I_k=(\frac{d}{dx})^k.$$
Thus (6.11) equals $`(K)`$ where $``$ is the fractional integral operator
$$:=C_1I_{3/2}x^{\frac{3}{2}}+C_2I_{\frac{1}{2}}x^{\frac{1}{2}}+C_3I_{\frac{1}{2}}x^{\frac{1}{2}}.$$
$`(6.13)`$
To solve for $`K`$ we apply $`I_{\frac{1}{2}}`$ to $`K`$ to get
$$C_1^{}\frac{d}{dx}^2(x^{\frac{3}{2}}K(x))+C_2^{}\frac{d}{dx}(x^{\frac{1}{2}}K)+C_3^{}x^{\frac{1}{2}}K=I_{\frac{1}{2}}K.$$
$`(6.14)`$
This equation determines $`K`$ up to a solution $`f`$ of the associated homogeneous equation, essentially an Euler equation. But also $`K=0`$ on $`[0,a(r_o)^2]`$ Since no homogeneous solution can have this property, $`K`$ is uniquely determined by this boundary condition.
It follows that both (6.9) and (6.12) are spectral invariants. But from $`a+b`$ and $`\frac{1}{a}+\frac{1}{b}`$ one can determine the pair $`(a,b)`$. Hence $`a^{}(r)`$ is determined from the spectrum. Since $`a(0)=0`$ this determines $`a`$ and hence the surface.
## 7 Appendix
The purpose of this appendix is to give an algorithm for calculating the higher order terms in the quasi-classical approximation of eigenvalues for 1 D Schrodinger operators $`\frac{h^2}{2}\frac{d^2}{dx^2}+V`$ with confining potentials. In particular, we carry out the calculation of the $`h^2E_n^{(2)}`$ term, which was used in the proof of the Final Lemma.
The algorithm is based on the Maslov method of canonical operators. Expositions and refinements of this method can be found, among other places, in Maslov’s book \[M\], in the article of Colin de Verdiere \[CV.3\] (and in its references), and in the recent book of Bates- Weinstein \[B.W\]. Although these references explain the construction of the canonical operator and prove the existence of complete quasi-classical eigenvalue expansions, they do not generally go on to describe the calculation of the terms. An exception is the original book of Maslov \[M\], which does calculate the first two or three terms; but the method of canonical operators is abandoned at this point in favor of some methods of special functions. As we will show, the canonical operator method gives the required corrections quite efficiently.
The set-up
We are concerned with the semi-classical eigenvalue problem:
$$\{\begin{array}{c}[\frac{h^2}{2}\frac{d^2}{dx^2}+V]\psi _n(x,h)=E_n(h)\psi _n(x,h)\hfill \\ \psi _n,\psi _m=\delta _{mn}+O(h^{\mathrm{}})\hfill \\ E_n(h)=E_n^{(1)}(h)+h^2E_n^{(2)}+h^3E_n^{(3)}+\mathrm{}\hfill \end{array}$$
The unknown function $`\psi _n(x,h)`$ is an oscillatory associated to a Lagrangean of the form $`\mathrm{\Lambda }_n:=\{H=E_n^{(1)}(h)\}`$ where $`E_n^{(1)}(h)`$ is determined by the Bohr-Sommerfeld-Maslov quantization condition:
$$\frac{1}{2\pi h}_{\mathrm{\Lambda }_n}\xi 𝑑x=n+\frac{1}{4}\mu .$$
Here, $`\mu `$ is the Maslov index of $`\mathrm{\Lambda }_n`$; it equals $`2`$ for connected level sets of Hamiltonians of the form $`H(x,\xi )=\xi ^2+V(x).`$
To find the higher order corrections $`E_n^{(k)}`$ we consider the Maslov canonical operator
$$𝒰_h:S^m(\mathrm{\Lambda }_n,\mathrm{\Omega }_{\frac{1}{2}})𝒪^m(\mathrm{I}\mathrm{R},\mathrm{\Lambda }_n).$$
We follow here the notation and terminology of \[B.W\]\[CV.3\]: $`S^m(\mathrm{\Lambda }_n,\mathrm{\Omega }_{\frac{1}{2}})`$ is the space of symbolic sections of the bundle of 1/2-densities times Maslov factors and $`𝒪^m(\mathrm{I}\mathrm{R},\mathrm{\Lambda }_n)`$ is the space of oscillatory integrals associated to $`\mathrm{\Lambda }_n`$. There is a natural symbol map in the reverse direction; any inverse modulo $`O(h)`$ is a quantization or canonical operator. Its existence is equivalent to the condition that $`\mathrm{\Lambda }_n`$ satisfy the BSM quantization condition. For background we again refer to \[CV.3\]\[B.W\].
We also denote by $`\mathrm{\Xi }_H`$ the Hamilton vector field of $`H`$, by $`_{\mathrm{\Xi }_H}`$ the Lie derivative on any bundle, by $`s`$ a nowhere vanishing section of $``$, and by $`\rho `$ a $`_{\mathrm{\Xi }_H}`$-invariant density on $`\mathrm{\Lambda }_n`$ (for a fixed $`n`$). By surjectivity of $`𝒰_h`$, we can write an oscillatory integral associated to $`\mathrm{\Lambda }_n`$ in the form
$$\psi _n(x,h)|dx|^{\frac{1}{2}}𝒰_h[e^{\frac{i}{h}\varphi }\underset{j=0}{\overset{\mathrm{}}{}}h^jf_j\rho ^{\frac{1}{2}}s].$$
Our aim is to determine the quasi-classical series $`E_n^{(j)}h^j`$ and coefficient functions $`f_j`$ for which the asymptotic eigenvalue problem is solvable.
We begin by constructing local solutions. Thus we first consider $`x`$-projectible pieces of $`\mathrm{\Lambda }_nT^{}\mathrm{I}\mathrm{R}`$: pieces which projects regularly from a neighborhood of $`\lambda \mathrm{\Lambda }_n`$ to a neighborhood of $`x\mathrm{I}\mathrm{R}.`$ Restricted to densities supported on such pieces, the Maslov canonical operator is truly canonical: if $`S(x)`$ is a phase locally parametrizing $`\mathrm{\Lambda }_n`$, then
$$𝒰_h[e^{\frac{i}{h}\varphi }\underset{j=0}{\overset{\mathrm{}}{}}h^jf_j\rho ^{\frac{1}{2}}s]=e^{\frac{i}{h}S}\underset{j=0}{\overset{\mathrm{}}{}}h^ja_j$$
for some smooth coefficients $`a_j`$. We may then substitute this expression into the eigenvalue problem and obtain eikonal and transport equations. The eikonal equation $`(S^{})^2+V(x)=E_n^{(1)}(h)=0`$ has been solved by our choice of phase, so the transport equations become:
$$\{\begin{array}{c}a_o\frac{d^2}{dx^2}S+2a_oS=0\hfill \\ a_1\frac{d^2}{dx^2}S+2a_1Si\frac{d^2}{dx^2}a_o=2iE_n^{(2)}a_o\hfill \\ a_2\frac{d^2}{dx^2}S+2a_2Si\frac{d^2}{dx^2}a_1=2i[E_n^{(3)}a_o+E_n^{(2)}a_1]\hfill \end{array}$$
and so on. As is well-known (cf. e.g. \[B.W\]), these equations may be put into geometric form by observing that $`S`$ is the projection to $`\mathrm{I}\mathrm{R}`$ of $`_{\mathrm{\Xi }_H}`$ and that $`[a\frac{d^2}{dx^2}S+2S]|dx|=div(a^2S)|dx|=_{\mathrm{\Xi }_H}(a^2|dx|).`$ Hence the transport equations become
$$\{\begin{array}{c}_{\mathrm{\Xi }_H}(a_o|dx|^{\frac{1}{2}})=0\hfill \\ _{\mathrm{\Xi }_H}(a_1|dx|^{\frac{1}{2}})=(2iE_n^{(2)}a_o+i\frac{d^2}{dx^2}a_o)|dx|^{\frac{1}{2}}\hfill \\ _{\mathrm{\Xi }_H}(a_2|dx|^{\frac{1}{2}})=2i[(E_n^{(3)}a_o+E_n^{(2)}a_1)+i\frac{d^2}{dx^2}a_1)]|dx|^{\frac{1}{2}}\hfill \end{array}$$
The invariant 1/2-density is given by the well-known formula $`a_o=(EV)^{\frac{1}{4}}`$, or equivalently by $`\frac{|dx|^{\frac{1}{2}}}{p^{\frac{1}{2}}}`$ on $`\mathrm{\Lambda }_n.`$ If we write $`a_j|dx|^{\frac{1}{2}}=f_j\rho ^{\frac{1}{2}}`$ then $`f_0=1,f_j=\frac{a_j}{a_o}`$ and the transport equations become
$$\{\begin{array}{c}\mathrm{\Xi }_Hf_1=(2iE_n^{(2)}+ia_o^1\frac{d^2}{dx^2}a_o)\hfill \\ \mathrm{\Xi }_Hf_2=2i[(E_n^{(3)}+E_n^{(2)}\frac{a_1}{a_o})+ia_o^1\frac{d^2}{dx^2}a_1)]\hfill \end{array}$$
Here, the expressions in $`a_o,a_1,\mathrm{}`$ are understood to have been lifted up to $`\mathrm{\Lambda }_n`$.
The eigenvalue corrections $`E_n^{(k)}`$ are determined by integrating both sides of the transport equations over the level $`\{H=E_n^{(1)}\}`$. Since the equations are solvable and since the left hand sides will integrate to zero, we get
$$\{\begin{array}{c}E_n^{(2)}=\frac{1}{2iT(E_n^{(1)})}_{\{H=E_n^{(1)}\}}[a_o^1\frac{d^2}{dx^2}a_o]\rho \hfill \\ E_n^{(3)}=\frac{1}{2iT(E_n^{(1)})}_{\{H=E_n^{(1)}\}}[E_n^{(2)}\frac{a_1}{a_o}+ia_o^1\frac{d^2}{dx^2}a_1]\rho \hfill \end{array}$$
Here, $`T(E)`$ is the period of $`\mathrm{\Xi }_H`$ at level $`E`$. Parametrizing $`\{H=E_n^{(1)}\}`$ as a graph over the $`x`$-axis away from the turning points, the invariant density takes the form $`\frac{dx}{\sqrt{EV}}`$ with $`E=E_1^{(1)}`$. Hence at least formally the eigenvalue corrections are given by
$$E_n^{(2)}=\frac{1}{2iT(E)}_{x_{}(E)}^{x_+(E)}[a_o^1\frac{d^2}{dx^2}a_o]\frac{dx}{\sqrt{EV}}=\frac{1}{2iT(E)}_{x_{}(E)}^{x_+(E)}(EV(x))^{\frac{1}{4}}\frac{d^2}{dx^2}(EV(x))^{\frac{1}{4}}𝑑x.$$
Unfortunately the integral is ill-defined due to the singularities at the turning points. The problem is that the Maslov operator cannot be defined near these points as a simple pull-back operator. Rather it should be defined as the composition of the Fourier transform with the pull-back operator defined over the $`\xi `$-projection. This problem and its solution constitute a key aspect of the Maslov method (in one dimension); we refer to \[CV.3\]\[B.W\] for extended discussions
The point which is important for us is that the Maslov method gives a regularization of the divergent integral. It works in the following way: we introduce a cut-off $`\psi _\delta `$ supported away from a $`\delta `$-neighborhood of the turning points $`x_\pm (E).`$ More precisely we define $`\psi _\delta ^\pm `$ on $`\mathrm{\Lambda }_n`$, equal to one on $`(2\delta ,\frac{1}{2}T(E)2\delta )`$ resp. $`(\frac{1}{2}T(E)+2\delta ,T(E)2\delta )`$ and equal to zero on $`(T(E)\delta ,\delta )`$ resp. $`(\frac{1}{2}T(E)\delta ,\frac{1}{2}T(E)+\delta ).`$ We then put
$$𝒰_h(f\rho se^{\frac{i}{h}\varphi }):=I_h(\psi _\delta f\rho se^{\frac{i}{h}\varphi })+J_h((1\psi _\delta )f\rho se^{\frac{i}{h}\varphi })$$
where $`I_h`$ is the pull-back to $`\mathrm{I}\mathrm{R}`$ under the phase parametrization by $`\xi =S^{}(x)`$ and where $`J_h`$ is the Fourier transform of the $`\xi `$-parametrization. The notation $`\psi _\delta `$ stands for $`\psi _\delta ^\pm `$. For details on $`I_h,J_h`$, see \[B.W\].
Returning to the previous calculation of eigenvalue corrections, we see that what is missing is the cut-off $`\psi _\delta `$ in the integrals and the contributions from $`J_h`$. We wish to avoid confronting the latter. Fortunately, it is not necessary to do so: the fact that the eigenvalues are independent of $`\delta `$ allows us to determine the $`J_h`$ (i.e. the turning point) contribution indirectly. To see this, we substitute the cut-off into the formula for $`E_n^{(2)}`$ to get
$$E_n^{(2)}=\frac{1}{2iT(E)}_{x_{}(E)}^{x_+(E)}(EV(x))^{\frac{1}{4}}\frac{d^2}{dx^2}[(EV(x))^{\frac{1}{4}}\psi _\delta ]𝑑x+II_\delta $$
with $`II_\delta `$ the contribution from $`J_h`$. Since the integral is now nicely convergent we can integrate by parts and simplify to the form
$$\frac{1}{16T(E)}_{x_{}(E)}^{x_+(E)}\frac{V^{}(x)^2}{(EV)^{\frac{5}{2}}}\psi _\delta 𝑑x\frac{1}{T(E)}_{x_{}(E)}^{x_+(E)}\frac{V^{}(x)}{(EV)^{\frac{3}{2}}}\psi _\delta ^{}(x)𝑑x.$$
The first term tends to
$$\frac{1}{12T(E)}\frac{d^2}{dE^2}_{x_{}(E)}^{x_+(E)}\frac{V^{}(x)^2}{(EV)^{\frac{1}{2}}}\psi _\delta 𝑑x$$
as $`\delta 0`$. The second expression only involves the Taylor expansion of $`V`$ near the turning points. Its singular part must be cancelled by the singular part of $`II_\delta `$, leaving a possible ‘residue’ as $`\delta 0.`$ We claim that this residue is zero: in fact, this is known to be the case since the first term is well-defined, independent of $`\delta `$, and agrees with the formula given in \[M\]. To give an independent proof that it vanishes, without analysing the $`J_h`$-term in detail, we observe that the limit contribution involves only the 2-jet of $`V`$ at the turning points. Hence it must agree with the corresponding expression for a harmonic oscillator at its turning points. But no such correction occurs. |
warning/0002/math0002251.html | ar5iv | text | # Higher homotopy groups of complements of complex hyperplane arrangements
## 1. Introduction
### 1.1. Background
One of the fundamental problems in the topological study of polynomial functions, $`f:(^{\mathrm{}},\mathrm{𝟎})(,0)`$, is the computation of the homotopy groups of the complement to the hypersurface $`V(f)=f^1(0)`$. A well-known algorithm for finding a finite presentation for $`\pi _1(^{\mathrm{}}V(f))`$ was given by Zariski and VanKampen in the early 1930’s. Much less is known about the higher homotopy groups of the complement, except when $`V(f)`$ is irreducible, in which case the Zariski-VanKampen method can be extended to give information about $`\pi _k(^{\mathrm{}}V(f))`$, for $`k>1`$, see .
In this paper, we concentrate on the simplest kind of polynomial $`f`$ for which the hypersurface $`V(f)`$ is not irreducible. Namely, suppose $`f`$ factors completely into distinct, degree one factors. Then $`f`$ is the defining polynomial of a hyperplane arrangement, $`𝒜`$, with union $`V(f)=_{H𝒜}H`$, and complement $`X(𝒜)=^{\mathrm{}}_{H𝒜}H`$. The cohomology ring of $`X=X(𝒜)`$ was computed by Brieskorn . Orlik and Solomon expressed $`H^{}(X)`$ in terms of the combinatorics of $`𝒜`$, encoded in the intersection lattice, $`L(𝒜)`$. In particular, the Poincaré polynomial, $`P_𝒜(T)=_{k=1}^{\mathrm{}}b_k(X)T^k`$, admits a simple combinatorial expression, see Orlik and Terao . On the other hand, the fundamental group of the complement, $`\pi _1(X)`$, is not determined by $`L(𝒜)`$ alone, as the example of Rybnikov shows.
For certain arrangements, all the higher homotopy groups of the complement vanish. Examples of such $`K(\pi ,1)`$ arrangements include the simplicial arrangements (Deligne ), and the supersolvable arrangements (Terao ). Examples of non-$`K(\pi ,1)`$ arrangements, and methods for detecting the first non-vanishing higher homotopy group of their complements, were given by Falk and Randell (see also the recent survey ).
The first (and, up to now, only) explicit computation of non-trivial higher homotopy groups of arrangement complements was made by Hattori . An arrangement $`𝒜`$ in $`^{\mathrm{}}`$, $`\mathrm{}>1`$, is called generic if, for all $`𝒜`$, the intersection $`_HH`$ has codimension $`\left|\right|`$ when $`\left|\right|\mathrm{}`$, and is empty when $`\left|\right|>\mathrm{}`$. The standard example is the Boolean arrangement of coordinate hyperplanes in $`^n`$, with complement $`(^{})^n`$. Hattori used the minimal cell decomposition of $`(^{})^nT^n`$ to find an explicit, minimal cell decomposition for the complement of an arbitrary generic arrangement. More precisely, if $`𝒜`$ is an arrangement of $`n`$ hyperplanes in general position in $`^{\mathrm{}}`$ ($`n>\mathrm{}`$), then $`X(𝒜)(T^n)^{(\mathrm{})}`$. From this decomposition, Hattori deduced:
1. $`\pi _1(X)=^n`$.
2. $`\pi _k(X)=0`$ for $`1<k<\mathrm{}`$.
3. $`\pi _{\mathrm{}}(X)`$ admits a free $`\pi _1`$-resolution of length $`n\mathrm{}`$.
The simplest example is that of $`3`$ generic affine lines in $`^2`$. In that case, the complement $`X`$ has the homotopy type of the $`2`$-skeleton of the $`3`$-torus, $`T^3=K(^3,1)`$. Looking at the universal cover, $`\stackrel{~}{X}`$, we thus see that $`\pi _2(X)=H_2(\stackrel{~}{X})`$ is a free $`^3`$-module, generated by the boundary of a cubical $`3`$-cell from the standard decomposition of $`\stackrel{~}{T^3}=^3`$.
### 1.2. Results
In this paper, we set out to generalize Hattori’s results to the wider class of hypersolvable arrangements. This combinatorially defined class, introduced in , includes both supersolvable arrangements and (cones of) generic arrangements.
A hypersolvable arrangement $`𝒜`$ admits a “supersolvable deformation,” $`\widehat{𝒜}`$, which preserves the collinearity relations. For example, if we start with $`n3`$ generic lines in $`^2`$, and take $`𝒜`$ to be the respective central arrangement of planes in $`^3`$, then $`\widehat{𝒜}`$ is the Boolean arrangement in $`^n`$. In general, $`X(𝒜)`$ has the same fundamental group as $`X(\widehat{𝒜})`$; see . Moreover, $`\pi _1(X(\widehat{𝒜}))`$ is a (special kind of) iterated semidirect product of finitely-generated free groups; see and Theorem 4.8. These facts together provide the generalization of Hattori’s result (A) to hypersolvable arrangements.
The key tool for generalizing (B) and (C) to complements of hypersolvable arrangements is the existence of minimal cell structures, on both $`X(𝒜)`$ and $`X(\widehat{𝒜})=K(\pi _1(X(𝒜)),1)`$.
To find a presentation for the first higher non-vanishing homotopy group, we thus turn to a general study of minimal cell decompositions. The idea is to get higher homotopy information on a connected, finite-type, CW-space $`X`$, by comparing it to its classifying space $`K(\pi ,1)`$, where $`\pi =\pi _1(X)`$. We are thus led to introduce a homotopy-type invariant of $`X`$, called the order of $`\pi _1`$-connectivity, which measures the rational-homology deviation of $`X`$ from asphericity:
$$p(X):=sup\{qb_r(X)=b_r(K(\pi ,1)),rq\}.$$
(If $`X`$ is $`1`$-connected, and $`H_{}(X)`$ is torsion-free, then $`p(X)`$ is the usual order of connectivity of $`X`$.)
A (connected, finite-type) CW-space $`X`$ is said to be minimal if it has a CW-structure with $`b_k`$ $`k`$-cells, for all $`k`$, where $`b_k`$ is the $`k`$-th Betti number of $`X`$. In §2, we will prove the following result.
###### Theorem 1.3.
Let $`X`$ and $`Y`$ be two minimal CW-complexes, with cohomology rings generated in degree $`1`$. Set $`\pi =\pi _1(X)`$ and $`p=p(X)`$. Assume that $`Y`$ is aspherical and that there is a cellular map, $`j:XY`$, such that $`j|{}_{X^{(p)}}{}^{}=\mathrm{id}`$. Then $`j`$ is a classifying map. Moreover:
1. $`X`$ is aspherical if and only if $`p=\mathrm{}`$.
2. If $`p<\mathrm{}`$, then the first non-trivial higher homotopy group of $`X`$ is $`\pi _p(X)`$, which has the following finite $`\pi `$-presentation:
(1.1)
$$(H_{p+2}(Y)H_{p+1}(X))\pi \stackrel{D_p:=_{p+2}+\stackrel{~}{j}_{p+1}}{}H_{p+1}(Y)\pi \pi _p(X)0,$$
where $`_{}`$ is the differential of the $`\pi `$-equivariant cellular chain complex of the universal cover of $`Y`$, and $`\stackrel{~}{j}_{}`$ is the $`\pi `$-chain map induced by the lift of $`j`$ to universal covers.
3. If $`p<\mathrm{}`$, then the group of $`\pi `$-coinvariants of $`\pi _p(X)`$ is isomorphic to $`H_{p+1}(Y,X)`$, where $`H_{}(Y,X):=\mathrm{coker}(j_{}:H_{}(X)H_{}(Y))`$.
The theorem provides a complete generalization of Hattori’s result (B) in this setting, and a partial generalization of (C). If $`dimX=p(X)`$ and $`Y`$ is a finite complex, the presentation (1.1) extends to a finite-length, free $`\pi `$-resolution of $`\pi _p(X)`$, see Remark 2.12. In particular, if $`X`$ is the complement of a generic arrangement, then $`p(X)=\mathrm{}`$, and (1.1) may be continued to Hattori’s resolution (C).
In §3, we follow a standard approach and extract from the above presentation of $`\pi _p(X)`$ more manageable invariants of homotopy type: the subvarieties of the complex torus $`(^{})^n`$, $`n=b_1(X)`$, defined by the Fitting ideals of $`\pi _p(X)_\pi ^n`$. In turn, we identify these varieties with the jumping loci for homology with coefficients in rank $`1`$ local systems of the pair $`(K(\pi ,1),X)`$.
We now return to the case where $`X=X(𝒜)`$ is the complement of an (essential) arrangement $`𝒜`$. From recent work of Dimca and Randell , we know that $`X`$ is minimal. We also know (from ) that $`H^{}(X)`$ is generated in degree $`1`$. Since $`X`$ may fail to possess any finite-type $`K(\pi ,1)`$, our approach does not work in this generality. If $`𝒜`$ is hypersolvable, though, we may take $`K(\pi ,1)=X(\widehat{𝒜})`$, where $`\widehat{𝒜}`$ is the supersolvable deformation of $`𝒜`$.
We devote §4 to the description of the first non-vanishing higher homotopy group of a hypersolvable (non-supersolvable) arrangement, $`𝒜`$. The combinatorics and the homotopy theory of the complement of $`𝒜`$ do not change, when passing to the associated essential arrangement, $`𝒜_{\mathrm{ess}}`$; see . Let $`\widehat{𝒜}`$ be the supersolvable deformation of the essential hypersolvable arrangement $`𝒜_{\mathrm{ess}}`$. Exploiting the method from \[9, Section 5\], we show in Theorems 4.2 and 4.8 (6) how to replace, up to homotopy, the inclusion $`J:X(𝒜_{\mathrm{ess}})X(\widehat{𝒜})`$ by a map, $`j:XY`$, which satisfies all hypotheses from Theorem 1.3. This leads to the following result.
###### Theorem 1.4.
Let $`𝒜`$ be a hypersolvable arrangement, with complement $`X=X(𝒜)`$, fundamental group $`\pi =\pi _1(X)`$, and order of $`\pi _1`$-connectivity $`p=p(X)`$. Then:
1. $`X`$ is aspherical $``$ $`𝒜`$ is supersolvable $``$ $`p=\mathrm{}`$.
2. If $`p<\mathrm{}`$, then the first non-vanishing higher homotopy group of $`X`$ is $`\pi _p(X)`$, with finite $`\pi `$-presentation (1.1).
3. If $`p<\mathrm{}`$, then the group of $`\pi `$-coinvariants of $`\pi _p(X)`$ is free abelian, with (strictly positive) combinatorially determined rank.
In Corollary 4.10(1), we give a combinatorial interpretation for $`p(X)`$. The precise formula for the coinvariants is provided by Theorem 4.12(3). A similar formula was obtained by Randell , for generic sections of aspherical arrangements—a class of arrangements which overlaps with the hypersolvable class, but neither includes it, nor is included in it.
### 1.5. Applications
Particularly simple is the case of affine line arrangements in $`^2`$. These arrangements represent both the simplest case of non-irreducible plane algebraic curves, and the simplest case of hyperplane arrangements. As such, they have been the object of intense investigation, see e.g. . For one, the fundamental group of an arbitrary hyperplane arrangement complement can be identified with the fundamental group of an affine line arrangement (by the Hamm-Lê theorem), thereby making line arrangements key to the understanding of all arrangements. For another, complements of affine line arrangements need not be aspherical (unlike, say, complements of weighted-homogeneous plane curves, which always are), thereby making for a richer object of topological study.
In §5, we consider affine line arrangements whose cones are hypersolvable. In Theorem 5.4, we go further, providing a minimal, finite-length resolution for $`\pi _2`$, which completely generalizes Hattori’s resolution (C) in this context. We also obtain some finer information about the $`\pi _1`$-module structure of $`\pi _2`$: it is neither projective (except in a very special, combinatorially decidable case, when it is free, with rank combinatorially determined), nor nilpotent (except if it is trivial).
Another class of arrangements which can be fairly well understood from our point of view is that of (hypersolvable) graphic arrangements. In §6, we implement in this setting our method for higher homotopy computations. The class of arrangements associated to graphs without $`3`$-cycles provides a natural, rich supply of hypersolvable arrangements, which are neither supersolvable nor generic, and for which homotopy information may be extracted directly from the graph. As an illustration, we exhibit two graphic arrangements, whose complements have the same $`\pi _1`$, but non-isomorphic $`\pi _2`$’s (when viewed as $`\pi _1`$-modules).
This version of our paper corrects an oversight from .
###### Acknowledgments.
This research was supported by the Volkswagen-Stiftung (RiP-program at Oberwolfach). The authors wish to thank the Mathematisches Forschungsinstitut for the warm hospitality and excellent facilities provided. The second author was also supported by an RSDF grant from Northeastern University.
## 2. Minimal cell decompositions and homotopy groups
### 2.1. Minimal cell decompositions
Given a space $`X`$, consider the following conditions on its homotopy type:
1. $`X`$ is homotopy equivalent to a connected, finite-type CW-complex;
2. The integral homology groups $`H_{}(X)`$ are torsion-free;
3. The cup-product map $`:^{}H^1(X)H^{}(X)`$ is surjective.
These three conditions abstract some well-known topological properties of complements of complex hyperplane arrangements. Next, we delineate a class of spaces that satisfy condition (i) and a much stronger form of (ii).
###### Definition 2.2.
A space $`X`$ is called minimal if $`X`$ has the homotopy type of a connected, finite-type, CW-complex $`K`$ such that
(2.1)
$$\mathrm{\#}\{q\text{-cells in }K\}=\mathrm{rank}H_q(X;),\text{for all }q0.$$
This definition implies at once that all the (abelian) groups $`H_q(X)`$ are finitely-generated and torsion-free. Consequently, we may unambiguously speak about the Betti numbers of $`X`$, $`b_q(X)`$, without specifying the coefficients.
Let $`X`$ be a minimal space, and let $`C_{}(X)`$ be the cellular chain complex of $`X`$, corresponding to a minimal CW-decomposition. Let $`\pi =\pi _1(X)`$ be the fundamental group, $`\pi `$ its group ring, and $`ϵ:\pi `$ the augmentation map. Let $`\stackrel{~}{X}`$ be the universal cover of $`X`$, and let $`(C_{}(\stackrel{~}{X}),d_{})`$ be the $`\pi `$-equivariant chain complex of $`\stackrel{~}{X}`$, with $`C_q(\stackrel{~}{X})=C_q(X)\pi `$ and $`d_q:C_q(\stackrel{~}{X})C_{q1}(\stackrel{~}{X})`$. (Note that, when turning left $`\pi `$-modules into right $`\pi `$-modules, one has to replace the action of $`x\pi `$ by that of $`x^1`$.) By minimality, all the boundary maps $`d_q`$ are $`ϵ`$-minimal, in the sense that $`d_q_\pi =0`$.
###### Example 2.3.
The standard example of a space admitting a minimal cell decomposition is the $`n`$-torus, $`T^n`$. Identifying $`\pi _1(T^n)=^n`$, with basis $`\{x_i\}_i`$, and $`C_q(T^n)=^q^n`$, with basis $`\{\sigma _I=\sigma _{i_1}\mathrm{}\sigma _{i_q}\}_I`$, the boundary map $`d_q:^q^n^n^{q1}^n^n`$ is given by $`d_q(\sigma _I)=_{r=1}^q(1)^{r1}\sigma _{I\{i_r\}}(x_{i_r}^11)`$.
###### Example 2.4.
More generally, let $`\pi =F_{d_n}_{\rho _{n1}}F_{d_{n1}}\mathrm{}_{\rho _1}F_{d_1}`$ be an iterated semidirect product of free groups, with $`\rho _i`$ acting as the identity in homology, and $`X=K(\pi ,1)`$ a corresponding Eilenberg-MacLane space. A finite, minimal cell decomposition of $`X`$ is given in : The number of cells is read off the Poincaré polynomial, $`P_X(T)=_{i=1}^n(1+d_iT)`$, and the ($`ϵ`$-minimal) boundary maps, $`d_q:C_q(\stackrel{~}{X})C_{q1}(\stackrel{~}{X})`$, are given explicitly in terms of Fox Jacobians of the monodromy operators $`\rho _i`$, see \[6, Thm. 2.10, Prop. 3.3, and Cor. 3.4\].
###### Remark 2.5.
Not all manifolds admit minimal cell decompositions. For example, if $`X`$ is the complement of a non-trivial knot in $`S^3`$, then $`X`$ has no minimal cell decomposition, not even up to $`q=1`$. See also the monograph by Sharko for various other definitions of minimality in related contexts.
Now assume $`X`$ is a minimal space for which there exists a minimal Eilenberg-MacLane space $`Y=K(\pi ,1)`$. Let $`j:XY`$ be a classifying map. Without loss of generality, we may assume $`j`$ respects the given (minimal) CW-decompositions on $`X`$ and $`Y`$. Then the chain map $`j_{}:C_{}(X)C_{}(Y)`$ lifts to an equivariant chain map $`\stackrel{~}{j}_{}:(C_{}(\stackrel{~}{X}),d_{})(C_{}(\stackrel{~}{Y}),_{})`$,
$$\stackrel{~}{j}_{}=\{\stackrel{~}{j}_q:H_q(X)\pi H_q(Y)\pi \}_{q0}.$$
### 2.6. Homotopy groups
We now analyze the homotopy groups of certain minimal spaces. In order to state our results, we need to introduce one more notion.
###### Definition 2.7.
Let $`X`$ be a space satisfying condition (i). Define the order of $`\pi _1`$-connectivity of $`X`$ to be
$$p(X):=sup\{qb_r(X;)=b_r(K(\pi _1(X),1);),rq\}.$$
###### Remark 2.8.
The terminology is borrowed from the simply-connected case: if $`\pi _1(X)=0`$ and $`X`$ also satisfies (ii), then $`p(X)`$ is the usual order of connectivity of $`X`$. Note that $`p(X)`$ is a positive integer, depending only on the homotopy type of $`X`$. Furthermore, if $`X`$ satisfies conditions (i)–(iii), then $`p(X)2`$.
###### Remark 2.9.
Set $`Y=K(\pi _1(X),1)`$, and consider a classifying map, $`j:XY`$. Assume both $`X`$ and $`Y`$ satisfy conditions (i)–(iii) from 2.1. These conditions readily imply that $`j`$ induces a split surjection on cohomology, and a split injection on homology. Consequently, $`j_r:H_r(X)H_r(Y)`$ is an isomorphism, for all $`rp(X)`$, and the groups $`H_{}(Y,X):=\mathrm{coker}(j_{}:H_{}(X)H_{}(Y))`$ fit into split exact sequences
(2.2)
$$0H_{}(X)\stackrel{j_{}}{}H_{}(Y)\stackrel{\mathrm{\Pi }_{}}{}H_{}(Y,X)0.$$
The next two results provide a complete proof of Theorem 1.3 from the Introduction.
###### Theorem 2.10.
Let $`X`$ and $`Y`$ be two minimal CW-complexes satisfying condition (iii) from §2.1. Set $`\pi =\pi _1(X)`$ and $`p=p(X)`$. Assume that $`Y`$ is aspherical, and that there is a cellular map, $`j:XY`$, such that the restriction of $`j`$ to the $`p`$-skeleton, $`X^{(p)}`$, is the identity. Then $`j`$ is a classifying map, and:
1. $`\stackrel{~}{X}`$ is $`(p1)`$-connected.
2. If $`p<\mathrm{}`$, then
(2.3)
$$(H_{p+2}(Y)\pi )(H_{p+1}(X)\pi )\stackrel{D_p:=_{p+2}+\stackrel{~}{j}_{p+1}}{}H_{p+1}(Y)\pi \pi _p(X)0$$
is a finite presentation of $`\pi _p(X)`$ as $`\pi `$-module.
###### Proof.
Since $`p2`$, $`Y`$ is a $`K(\pi ,1)`$ and $`j`$ is a classifying map.
(1) We have to show that $`\pi _q(\stackrel{~}{X})=0`$, for $`q<p`$. Of course, $`\pi _1(\stackrel{~}{X})=0`$. Fix $`1<q<p`$, and assume that $`\pi _r(\stackrel{~}{X})=0`$, for $`r<q`$. By the Hurewicz isomorphism theorem, $`\pi _q(\stackrel{~}{X})=H_q(\stackrel{~}{X})`$. By minimality of $`X`$ and $`Y`$, we have a commuting ladder between the (equivariant) chain complexes of $`\stackrel{~}{X}`$ and $`\stackrel{~}{Y}`$:
$$\begin{array}{cccccc}C_{}(\stackrel{~}{X}):& H_{q+1}(X)\pi & \stackrel{d_{q+1}}{}& H_q(X)\pi & \stackrel{d_q}{}& H_{q1}(X)\pi \\ & \stackrel{~}{j}_{q+1}& & \stackrel{~}{j}_q& & \stackrel{~}{j}_{q1}& & \\ C_{}(\stackrel{~}{Y}):& H_{q+1}(Y)\pi & \stackrel{_{q+1}}{}& H_q(Y)\pi & \stackrel{_q}{}& H_{q1}(Y)\pi \end{array}$$
The three vertical arrows are isomorphisms, since $`\stackrel{~}{j}_r=\mathrm{id}`$, for $`rp(X)`$. It follows that $`H_q(\stackrel{~}{X})=H_q(\stackrel{~}{Y})`$. But $`\stackrel{~}{Y}`$ is acyclic, and so $`\pi _q(\stackrel{~}{X})=0`$.
(2) Consider the commuting diagram
$$\begin{array}{ccccccc}& & H_{p+1}(X)\pi & \stackrel{d_{p+1}}{}& H_p(X)\pi & \stackrel{d_p}{}& H_{p1}(X)\pi \\ & & \stackrel{~}{j}_{p+1}& & \mathrm{id}& & \mathrm{id}& & \\ H_{p+2}(Y)\pi & \stackrel{_{p+2}}{}& H_{p+1}(Y)\pi & \stackrel{_{p+1}}{}& H_p(Y)\pi & \stackrel{_p}{}& H_{p1}(Y)\pi \end{array}$$
By Part (1) and Hurewicz, $`\pi _p(X)=H_p(\stackrel{~}{X})`$. A diagram chase yields isomorphisms
$$H_p(\stackrel{~}{X})=\frac{\mathrm{ker}d_p}{\mathrm{im}d_{p+1}}=\frac{\mathrm{ker}_p}{\mathrm{im}d_{p+1}}=\frac{\mathrm{im}_{p+1}}{\mathrm{im}(_{p+1}\stackrel{~}{j}_{p+1})}\stackrel{_{p+1}}{}\frac{H_{p+1}(Y)\pi }{\mathrm{im}_{p+2}+\mathrm{im}\stackrel{~}{j}_{p+1}}=\mathrm{coker}(D_p),$$
and we are done. ∎
###### Corollary 2.11.
With assumptions as above, and if $`p=p(X)<\mathrm{}`$, then the group of coinvariants of $`\pi _p(X)`$ under the action of $`\pi =\pi _1(X)`$ is given by
$$(\pi _p(X))_\pi =H_{p+1}(Y,X).$$
In particular, $`(\pi _p(X))_\pi 0`$.
###### Proof.
From presentation (2.3), we compute $`(\pi _p(X))_\pi =\mathrm{coker}(D_p_\pi )`$. At the same time, $`\mathrm{coker}(D_p_\pi )=\mathrm{coker}(\stackrel{~}{j}_{p+1}_\pi )`$, since $`_{p+2}`$ is $`ϵ`$-minimal, and $`\stackrel{~}{j}_{p+1}_\pi =j_{p+1}`$, since $`\stackrel{~}{j}_{}`$ lifts $`j_{}`$. This proves the first claim. To finish the proof, note that $`H_{p+1}(Y,X)0`$, by the definition of $`p`$. ∎
Assuming $`p<\mathrm{}`$ in Theorem 2.10, one may also consider the following $`\pi `$-linear map:
(2.4)
$$H_{p+2}(Y)\pi \stackrel{\mathrm{\Delta }_p:=(\mathrm{\Pi }_{p+1}\mathrm{id})_{p+2}}{}H_{p+1}(Y,X)\pi .$$
In the next two remarks, we present some cases when the presentation matrix $`D_p`$ from (2.3) may be replaced by the (simpler, $`ϵ`$-minimal) presentation matrix $`\mathrm{\Delta }_p`$ from (2.4).
###### Remark 2.12.
If $`X`$ has the homotopy type of a CW-complex of dimension $`p=p(X)`$, and $`Y`$ is finite, then the presentation (2.3) for $`\pi _p(X)`$ may be continued to a free $`\pi `$-resolution of length $`dp`$, where $`d=dimY`$. (We shall encounter such a situation later on, in Theorem 5.4.) Indeed, let $`(C_{}(\stackrel{~}{Y}),_{})`$ be the $`\pi `$-equivariant cellular chain complex of $`\stackrel{~}{Y}`$. Note that $`H_{p+1}(X)=0`$ (since $`dimX=p`$), and so $`H_{p+1}(Y,X)=H_{p+1}(Y)`$ and $`\mathrm{\Delta }_p=_{p+2}=D_p`$. Hence, $`\pi _p(X)`$ has finite, free, $`ϵ`$-minimal, resolution
$$0C_d(\stackrel{~}{Y})\stackrel{_d}{}C_{d1}(\stackrel{~}{Y})\mathrm{}C_{p+2}(\stackrel{~}{Y})\stackrel{_{p+2}}{}C_{p+1}(\stackrel{~}{Y})\pi _p(X)0.$$
###### Remark 2.13.
An especially simple situation where Theorem 2.10 applies is as follows. Let $`Y`$ be a minimal $`K(\pi ,1)`$-complex satisfying condition (iii), and let $`XY`$ be a proper, connected subcomplex, such that $`X^{(2)}=Y^{(2)}`$. Since $`Y`$ is minimal, $`X`$ is also minimal, and (iii) also holds for $`X`$. Since $`X`$ and $`Y`$ share the same $`2`$-skeleton, the inclusion $`j:XY`$ is a classifying map, inducing an isomorphism on $`\pi _1`$. Since $`X`$ is a subcomplex of $`Y`$, we have an exact sequence of cellular chain complexes, $`0C_{}(X)\stackrel{j_{}}{}C_{}(Y)\stackrel{\mathrm{pr}_{}}{}C_{}(Y,X)0`$, and $`\mathrm{\Pi }_{}=\mathrm{pr}_{}:H_{}(Y)H_{}(Y,X)`$. Moreover,
$$p(X)=\mathrm{max}\{q\mathrm{\#}\{r\text{-cells of }X\}=\mathrm{\#}\{r\text{-cells of }Y\},rq\},$$
and $`p(X)<\mathrm{}`$. Therefore, $`\pi _k(\stackrel{~}{X})=0`$, for $`k<p=p(X)`$, and $`\pi _p(X)`$ has a finite $`\pi `$-presentation, given in (2.3). Obviously, $`\stackrel{~}{j}_{p+1}=j_{p+1}\mathrm{id}`$, since $`X^{(p+1)}`$ is a subcomplex of $`Y^{(p+1)}`$. Hence, $`\mathrm{\Pi }_{p+1}\mathrm{id}`$ induces an isomorphism between $`\mathrm{coker}(D_p)`$ and $`\mathrm{coker}(\mathrm{\Delta }_p)`$.
###### Remark 2.14.
Most of the results in this section have only relevance for non-simply-connected spaces. Indeed, if $`X`$ has the homotopy type of a finite-type CW-complex, and $`\pi _1(X)=0`$, then $`X`$ cannot satisfy condition (iii), unless $`X`$ is contractible.
On the other hand, if $`X`$ is $`1`$-connected and satisfies conditions (i) and (ii), then $`X`$ has the homotopy type of a minimal CW-complex $`K`$. The complex $`K`$ may be obtained from a bouquet of spheres $`_{b_{p+1}(X)}S^{p+1}`$, where $`p=p(X)`$, by attaching suitable cells. For details of the proof, see Anick , where the notion of minimality for simply-connected spaces was first introduced.
## 3. Fitting ideals and jumping loci
The $`\pi `$-module $`\pi _p(X)`$ determined in Theorem 2.10 can be rather intractable. We now associate to $`\pi _p(X)`$ a more manageable module (over a commutative ring), and extract from it invariants that can be understood as jumping loci for homology with coefficients in rank $`1`$ local systems.
### 3.1. Fitting ideals
Let $`\pi `$ be a group, with abelianization $`\pi ^{\mathrm{ab}}^n`$, and let $`M`$ be a finitely-presented module over $`\pi `$. Let $`\stackrel{~}{M}=M_\pi ^n`$ be the module over $`^n`$ obtained by extending scalars via $`\pi \stackrel{ab}{}^n`$. For $`k0`$, let $`F_k(\stackrel{~}{M})`$ be the corresponding $`k^{\text{th}}`$ Fitting ideal, generated by the codimension $`k1`$ minors of a presentation matrix for $`\stackrel{~}{M}`$. As is well-known, the Fitting ideals are independent of the choice of presentation, see e.g. \[10, p. 493\]. Now fix a basis $`\{x_1,\mathrm{},x_n\}`$ for $`^n`$. Then, the group ring $`^n=^n`$ may be identified with $`[x_1^{\pm 1},\mathrm{},x_n^{\pm 1}]`$, the coordinate ring of the complex algebraic torus $`(^{})^n`$. For each $`k0`$, the $`k^{\text{th}}`$ Fitting ideal of $`\stackrel{~}{M}`$ defines a subvariety of this torus,
$$V_k(M):=\{t(^{})^ng(t)=0,gF_k(\stackrel{~}{M})\}.$$
Alternatively, $`V_k(M)`$ can be described as the variety defined by the annihilator of $`^k\stackrel{~}{M}`$. Indeed, $`\mathrm{Rad}(F_k(\stackrel{~}{M}))=\mathrm{Rad}(\mathrm{ann}(^k\stackrel{~}{M}))`$, see \[10, pp. 511-513\].
Let $`X`$ be a path-connected space. Assume that $`\pi =\pi _1(X)`$ has abelianization $`^n`$, and that $`\pi _p(X)`$ is a finitely presented $`\pi `$-module.
###### Definition 3.2.
The $`k`$-th Fitting variety of $`\pi _p(X)`$ is the subvariety of the complex algebraic $`n`$-torus, $`V_k(\pi _p(X))`$, defined as above.
A standard argument (which we include for the sake of completeness) shows that the Fitting varieties are invariants of the $`\pi _1`$-module $`\pi _p`$.
###### Proposition 3.3.
Suppose $`f:\pi _1(X)\pi _1(X^{})`$ and $`g:\pi _p(X)\pi _p(X^{})`$ are compatible isomorphisms, i.e., $`g(xm)=f(x)g(m)`$, for all $`x\pi _1(X)`$ and $`m\pi _p(X)`$. Then there is a monomial isomorphism $`\mathrm{\Phi }:(^{})^n(^{})^n`$ such that $`\mathrm{\Phi }(V_k(\pi _p(X))=V_k(\pi _p(X^{}))`$.
###### Proof.
The extension of $`f_{\mathrm{ab}}:H_1(X)H_1(X^{})`$ to group rings gives rise to an isomorphism between $`\stackrel{~}{M}=\pi _p(X)_{\pi _1(X)}H_1(X)`$ and $`\stackrel{~}{M^{}}=\pi _p(X^{})_{\pi _1(X^{})}H_1(X^{})`$, and thus maps $`\mathrm{Rad}(\mathrm{ann}(^k\stackrel{~}{M}))`$ bijectively to $`\mathrm{Rad}(\mathrm{ann}(^k\stackrel{~}{M^{}}))`$.
Now identify $`H_1(X)`$ and $`H_1(X^{})`$ with $`^n`$, and let $`\varphi =\left(\varphi _{i,j}\right):^n^n`$ be the matrix of $`f_{\mathrm{ab}}`$ under this identification. Let $`\mathrm{\Phi }:(^{})^n(^{})^n`$ be the corresponding monomial isomorphism, given by $`\mathrm{\Phi }(t_i)=t_1^{\varphi _{i,1}}\mathrm{}t_n^{\varphi _{i,n}}`$. It is readily verified that $`\mathrm{\Phi }`$ preserves the Fitting varieties. ∎
###### Corollary 3.4.
For each $`k0`$, the monomial isomorphism type of $`V_k(\pi _p(X))`$ is a homotopy type invariant for $`X`$. ∎
Now let $`X`$ be a space satisfying the conditions of Theorem 2.10(2). Then $`\pi ^{\mathrm{ab}}^n`$, where $`n=b_1(X)`$. Let $`p=p(X)`$ be the order of $`\pi _1`$-connectivity of $`X`$. The Fitting ideals of $`\pi _p(X)_\pi ^n`$, and the varieties defined by them, may be computed from the presentation matrix $`D_p_\pi ^n`$, where $`D_p`$ is the matrix given in (2.3).
### 3.5. Characteristic varieties
There is another, well-known way to associate subvarieties of the complex algebraic torus to a space $`X`$ satisfying conditions (i) and (ii) from §2.1. For positive integers $`k`$ and $`p`$, set $`V_k^p(X)=\{t(^{})^ndim_{}H_p(C_{}(\stackrel{~}{X})_\pi _t)k\}`$, where $`n=b_1(X)`$, and $`_t`$ is the $`\pi `$-module $``$, given by the representation $`\pi \stackrel{\text{ab}}{}^n^{}`$, gotten by sending $`x_i`$ to $`t_i`$. This is an algebraic subvariety of $`(^{})^n`$, called the $`(p,k)`$-characteristic variety of $`X`$. It is straightforward to extend this definition to the relative setting, as follows.
###### Definition 3.6.
Let $`(Y,X)`$ be a CW-pair of spaces satisfying conditions (i) and (ii), and such that the inclusion $`XY`$ induces an isomorphism on $`\pi _1`$. Set $`n=b_1(X)`$. For $`k`$, $`p>0`$, the $`(p+1,k)`$-characteristic variety of $`(Y,X)`$ is the subvariety of the complex algebraic $`n`$-torus defined by
$$V_k^{p+1}(Y,X):=\{t(^{})^ndim_{}H_{p+1}(C_{}(\stackrel{~}{Y},\stackrel{~}{X})_\pi _t)k\}.$$
The characteristic variety $`V_k^1(X)`$ was interpreted by Hironaka as the variety defined by $`F_{k+1}(X)`$, the ideal generated by the codimension $`k`$ minors of the Alexander matrix of $`\pi _1(X)`$:
$$V_k^1(X)=V(F_{k+1}(X))$$
The next result provides a higher-dimensional analogue of Hironaka’s theorem, in a relative setting.
###### Proposition 3.7.
Let $`Y`$ be a minimal $`K(\pi ,1)`$-complex satisfying condition (iii) from §2.1, and let $`XY`$ be a proper, connected subcomplex, such that $`X^{(2)}=Y^{(2)}`$. Set $`p=p(X)`$. Then, for all $`k0`$,
$$V_k(\pi _p(X))=V_k^{p+1}(Y,X).$$
###### Proof.
By Remark 2.13, $`\pi _p(X)`$ is a finitely presented $`\pi `$-module, with presentation matrix $`\mathrm{\Delta }_p`$ given in (2.4). Hence, $`V_k(\pi _p(X))=\{t(^{})^ndim\mathrm{coker}\mathrm{\Delta }_p(t)k\}`$, where $`\mathrm{\Delta }_p(t)`$ is the evaluation of the matrix of Laurent polynomials $`\mathrm{\Delta }_p_\pi ^n`$ at $`x_i=t_i`$.
Let $`j:XY`$ be the inclusion. The lift to universal covers, $`\stackrel{~}{j}:\stackrel{~}{X}\stackrel{~}{Y}`$, gives rise to an exact sequence of $`\pi `$-equivariant chain complexes, a fragment of which is shown below:
(3.1)
Now fix $`t(^{})^n`$. Tensoring (3.1) over $`\pi `$ with $``$, via the representation $`t:\pi ^{}`$, yields the commuting diagram
(Note also that $`H_p(Y,X)=0`$, by the definition of $`p`$.) Chasing this diagram, we see that $`H_{p+1}(\stackrel{~}{Y},\stackrel{~}{X};_t)=\mathrm{coker}\overline{}_{p+2}(t)=\mathrm{coker}\mathrm{\Delta }_p(t)`$. From Definition 3.6, we have $`V_k^{p+1}(Y,X)=\{t(^{})^ndimH_{p+1}(\stackrel{~}{Y},\stackrel{~}{X};_t)k\}`$, and we are done. ∎
## 4. Higher homotopy groups of hypersolvable arrangements
We now apply the machinery developed in §2 to the class of spaces we had in mind all along: complements of complex hyperplane arrangements. We begin with a brief review of basic notions and relevant general results.
### 4.1. Minimal cell decompositions of arrangements
A (complex) hyperplane arrangement is a finite set, $`𝒜`$, of codimension-$`1`$ affine subspaces in a finite-dimensional complex vector space, $`V`$. The two main objects associated to an arrangement $`𝒜`$ are its complement, $`X(𝒜)=V_{H𝒜}H`$, and its intersection lattice, $`L(𝒜)=\{_HH𝒜\}`$. A general reference for the subject is the book by Orlik and Terao .
Since $`X(𝒜)`$ is the complement of a complex hypersurface, it has the homotopy type of a finite CW-complex, and thus satisfies condition (i) from §2.1. Explicit regular CW-complexes (of dimension equal to $`dim_{}V`$) onto which $`X(𝒜)`$ deform-retracts were given by Salvetti (in the complexified-real case), and by Björner and Ziegler (in the general case). Neither of these complexes, though, is minimal.
In a first preprint version of , we proved that the complements of arbitrary complex hyperplane arrangements satisfy the minimality condition (2.1), up to $`q=2`$, by combining results from and . Since then, the minimality question for arrangement complements, raised in preprint version $`1`$ of , has been solved in the affirmative by Dimca and Randell (independently). Using Morse theory, they proved the following.
###### Theorem (Dimca , Randell ).
Let $`𝒜`$ be a complex hyperplane arrangement, with complement $`X(𝒜)`$. Then $`X(𝒜)`$ is minimal, i.e., it admits a finite cell decomposition with number of $`q`$-cells equal to the $`q`$-th Betti number, for all $`q`$.
As noted in , complements of generic projective hypersurfaces fail to possess minimal cell structures. This indicates that minimality is a strong topological peculiarity of complements of complex arrangements.
Assume now that $`𝒜`$ is an essential $`k`$-generic section of an essential, central, aspherical arrangement $`\widehat{𝒜}`$, with $`k2`$. That is, $`\widehat{𝒜}`$ is an essential arrangement of hyperplanes passing through the origin of $`V`$, the complement $`X(\widehat{𝒜})`$ is aspherical, and there is a complex linear subspace, $`UV`$, such that $`U`$ is $`L_k(\widehat{𝒜})`$-generic (in the sense of definition §$`5(1)`$ from ), and the restriction $`\widehat{𝒜}^U`$ is equal to $`𝒜`$. In this framework, the refined Morse-theoretic analysis from \[9, Section 5\] leads to the following improved minimality result.
###### Theorem 4.2.
Let $`𝒜`$ be an essential, $`k`$-generic section ($`k2`$) of an essential, aspherical arrangement $`\widehat{𝒜}`$, as above. Set $`p=p(X(𝒜))`$. Then there is a cellular classifying map between minimal CW-complexes with cohomology rings generated in degree $`1`$, $`j:XY`$, which restricts to the identity on $`p`$-skeleta, and has the homotopy type of the inclusion, $`X(𝒜)X(\widehat{𝒜})`$.
###### Proof.
Denote by $`X^{}(𝒜)`$ and $`X^{}(\widehat{𝒜})`$ the complements of the associated projective arrangements in the projective spaces $`(U)`$ and $`(V)`$, respectively. The triviality of Hopf fibrations of arrangement complements readily gives homotopy equivalences, $`X(𝒜)X^{}(𝒜)\times S^1`$ and $`X(\widehat{𝒜})X^{}(\widehat{𝒜})\times S^1`$, and implies that $`p(X^{}(𝒜))=p`$.
Propositions $`14`$ and $`15`$ from may be used to replace, up to homotopy, the inclusion $`X^{}(𝒜)X^{}(\widehat{𝒜})`$ by a cellular map between minimal CW-complexes, $`j^{}:X^{}Y^{}`$, which restricts to the identity on $`p`$-skeleta; see the discussion preceding Theorem 16 from . Let $`X=X^{}\times S^1`$ and $`Y=Y^{}\times S^1`$. As shown by Brieskorn , all complements of complex hyperplane arrangements satisfy condition (iii) from §2.1. Thus, $`X`$ and $`Y`$ satisfy that condition, too. Now set $`j:=j^{}\times \mathrm{id}:XY`$. The claimed properties of the map $`j`$ follow from the corresponding properties of $`j^{}`$, and the fact that $`p2`$ (by Remark 2.8), which guarantees that $`j`$ is a classifying map. ∎
### 4.3. OS-algebras
The minimality result from Theorem 4.2 opens the way for using our approach to generalize Hattori’s results to a wider class of arrangements. But first, we need to recall an important result of Orlik and Solomon , which gives a combinatorial description of cohomology rings of arrangement complements.
Let $`𝒜=\{H_1,\mathrm{},H_n\}`$ be a central arrangement. By definition, the OS-algebra of $`𝒜`$ is
(4.1)
$$\mathrm{A}^{}(𝒜)=^{}(e_1,\mathrm{},e_n)/\left(e_{}|𝒜\text{and}\mathrm{codim}\underset{H}{}H<\left|\right|\right),$$
where $`^{}(e_1,\mathrm{},e_n)`$ is the exterior algebra over $``$ on generators $`e_i`$ in degree $`1`$, and for $`=\{H_{i_1},\mathrm{},H_{i_r}\}`$, $`e_{}=e_{i_1}\mathrm{}e_{i_r}`$ and $`e_{}=_q(1)^{q1}e_{i_1}\mathrm{}\widehat{e_{i_q}}\mathrm{}e_{i_r}`$. There is then an isomorphism of graded algebras,
(4.2)
$$\mathrm{OS}:H^{}(X(𝒜))\mathrm{A}^{}(𝒜).$$
Under this identification, the basis $`\{e_1,\mathrm{},e_n\}`$ of $`\mathrm{A}^1(𝒜)`$ is dual to the basis of $`H_1(X(𝒜))`$ given by the meridians of the hyperplanes, see . With respect to a fixed ordering of the hyperplanes, a canonical basis for $`\mathrm{A}^{}(𝒜)`$ is the no broken circuits (or, nbc) basis, see .
There is another, closely related, graded algebra, $`\overline{\mathrm{A}}^{}(𝒜)`$, called the quadratic Orlik-Solomon algebra, defined as the quotient of $`^{}(e_1,\mathrm{},e_n)`$ by relations of the form $`e_{}`$, for all $`𝒜`$ such that $`\mathrm{codim}_HH<\left|\right|`$ and $`\left|\right|=3`$, see . Clearly, the algebra $`\mathrm{A}^{}(𝒜)`$ is a quotient of $`\overline{\mathrm{A}}^{}(𝒜)`$, and the two algebras coincide up to degree $`2`$. Denote by
(4.3)
$$\pi _𝒜^{}:\overline{\mathrm{A}}^{}(𝒜)\mathrm{A}^{}(𝒜)$$
the canonical projection. Also denote by $`P_𝒜(T)`$ the Poincaré polynomial of $`\mathrm{A}^{}(𝒜)`$, and by $`\overline{P}_𝒜(T)`$ that of $`\overline{\mathrm{A}}^{}(𝒜)`$. It follows at once that $`\overline{P}_𝒜(T)P_𝒜(T)`$ (coefficientwise inequality).
### 4.4. Supersolvable and hypersolvable arrangements
Perhaps the best understood arrangements are the supersolvable (or, fiber-type) arrangements, introduced by Falk and Randell in . A central arrangement $`𝒜`$ is called supersolvable if its intersection lattice is supersolvable, in the sense of Stanley . For our purposes here, another (equivalent) combinatorial definition will be, however, more useful; see Definition 4.6. The standard example is the braid arrangement in $`^{\mathrm{}}`$, $`_{\mathrm{}}=\{\mathrm{ker}(z_iz_j)\}_{1i<j\mathrm{}}`$, with $`L(_{\mathrm{}})=𝒫_{\mathrm{}}`$, the partition lattice, and $`\pi _1(X(_{\mathrm{}}))=P_{\mathrm{}}`$, the pure braid group on $`\mathrm{}`$ strings. It follows from a theorem of Terao and results in that the complement of an arbitrary supersolvable arrangement is a $`K(\pi ,1)`$.
The class of hypersolvable arrangements actually motivated the framework for our Theorem 2.10. We start by reviewing the definition and basic properties of such arrangements.
Let $`𝒜=\{H_1,\mathrm{},H_n\}`$ be a central arrangement in the complex vector space $`V`$. Denote also by $`𝒜=\{\alpha _1,\mathrm{},\alpha _n\}(V^{})`$ its set of defining equations, viewed as points in the dual projective space. Let $`𝒜`$ be a proper, non-empty sub-arrangement, and set $`\overline{}:=𝒜`$. We say that $`(𝒜,)`$ is a solvable extension if the following conditions are satisfied (see ):
1. No point $`a\overline{}`$ sits on a projective line determined by $`\alpha ,\beta `$.
2. For every $`a,b\overline{}`$, $`ab`$, there exists a point $`\alpha `$ on the line passing through $`a`$ and $`b`$. (In the presence of condition (I), this point is uniquely determined, and will be denoted by $`f(a,b)`$.)
3. For every distinct points $`a,b,c\overline{}`$, the three points $`f(a,b)`$, $`f(a,c)`$, and $`f(b,c)`$ are either equal or collinear.
Note that only two possibilities may occur: either $`\mathrm{rank}(𝒜)=\mathrm{rank}()+1`$ (fibered case), or $`\mathrm{rank}(𝒜)=\mathrm{rank}()`$ (singular case); see \[21, Lemma 1.3(i)\].
###### Definition 4.5 ().
The arrangement $`𝒜`$ is called hypersolvable if it has a hypersolvable composition series, i.e., an ascending chain of sub-arrangements, $`𝒜_1\mathrm{}𝒜_i𝒜_{i+1}\mathrm{}𝒜_{\mathrm{}}=𝒜`$, where $`\mathrm{rank}𝒜_1=1`$, and each extension $`(𝒜_{i+1},𝒜_i)`$ is solvable.
The length of a composition series depends only on $`𝒜`$; it will be denoted by $`\mathrm{}(𝒜)`$. Note that the property of being hypersolvable is purely combinatorial. In fact, given an arrangement $`𝒜`$, one can decide whether it is hypersolvable or not, only from the elements of rank one and two of $`L(𝒜)`$, since the definitions only involve the collinearity relations in $`𝒜`$. The class of hypersolvable arrangements includes supersolvable arrangements, cones of generic arrangements (for which $`\mathrm{}(𝒜)=\left|𝒜\right|`$), and many others, see , and the examples in §§4.14, 5.6, and 6.8.
It is appropriate to mention here that one may assume from now on, whenever necessary, that a given hypersolvable arrangement $`𝒜`$ is also essential. Indeed, one knows how to associate to an arbitrary central arrangement $`𝒜`$ an essential arrangement, $`𝒜_{\mathrm{ess}}`$, without changing the homotopy types of the complements, $`X(𝒜)`$ and $`X^{}(𝒜)`$, and the intersection lattice $`L(𝒜)`$; see \[26, p. 197\].
The connection between hypersolvable and supersolvable (or fiber-type) arrangements comes from the following fact, which is implicit in \[21, Lemma 4.5\], and is explicitly proved in \[22, Prop. 1.3(i)\]. If the solvable extension $`(𝒜,)`$ is fibered, then there is a Serre fibration $`X(𝒜)X()`$, with homotopy fiber $`\{m\text{ points}\}`$, where $`m=|\overline{}|`$. It follows from \[22, Prop. 1.3\] that the (topological) definition of fiber-type arrangements may be rephrased in hypersolvable terms, as follows.
###### Definition 4.6 ().
The arrangement $`𝒜`$ is supersolvable (or, fiber-type) if it has a supersolvable composition series, that is, a hypersolvable composition series as in Definition 4.5, for which all extensions are fibered.
We thus see that all fiber-type arrangements are hypersolvable. On the other hand, one knows from \[21, Thm. D\] that a hypersolvable arrangement $`𝒜`$ is aspherical if and only if $`𝒜`$ is fiber-type, which happens precisely when $`\mathrm{}(𝒜)=\mathrm{rank}(𝒜)`$.
### 4.7. Supersolvable deformations
Our basic tool for the topological study of hypersolvable arrangements is the following theorem, which mostly puts together and organizes a number of known results.
###### Theorem 4.8.
Let $`𝒜`$ be an essential, hypersolvable arrangement, with composition series $`𝒜_1\mathrm{}𝒜_{\mathrm{}}=𝒜`$, and exponents $`d_i:=\left|𝒜_i𝒜_{i1}\right|`$. Set $`\pi =\pi _1(X(𝒜))`$ and $`p=p(X(𝒜))`$. Then $`𝒜`$ is a $`2`$-generic section of an essential, supersolvable arrangement $`\widehat{𝒜}`$, called the supersolvable deformation of $`𝒜`$, such that:
1. $`\widehat{𝒜}`$ has a supersolvable composition series, $`\widehat{𝒜}_1\mathrm{}\widehat{𝒜}_{\mathrm{}}=\widehat{𝒜}`$, with $`|\widehat{𝒜}_i|=\left|𝒜_i\right|`$, for $`1i\mathrm{}`$.
2. $`X(\widehat{𝒜})`$ sits atop a tower, $`X(\widehat{𝒜}_{\mathrm{}})\stackrel{p_{\mathrm{}}}{}X(\widehat{𝒜}_\mathrm{}1)\mathrm{}X(\widehat{𝒜}_2)\stackrel{p_2}{}X(\widehat{𝒜}_1)=^{}`$, of Serre fibrations, $`p_i:X(\widehat{𝒜}_i)X(\widehat{𝒜}_{i1})`$, with fiber $`\{d_i\text{ points}\}`$, and monodromy $`\rho _{i1}:\pi _1(X(\widehat{𝒜}_{i1}))P_{d_i}\mathrm{Aut}(F_{d_i})`$.
3. $`X(\widehat{𝒜})`$ is a $`K(\pi ,1)`$ space. The fundamental group admits an iterated semidirect product decomposition, $`\pi =F_d_{\mathrm{}}_{\rho _\mathrm{}1}\mathrm{}_{\rho _1}F_{d_1}`$.
4. There is a canonical isomorphism, $`\overline{\mathrm{OS}}:H^{}(X(\widehat{𝒜}))\overline{\mathrm{A}}^{}(𝒜)`$. Under this isomorphism, and the isomorphism $`\mathrm{OS}:H^{}(X(𝒜))\mathrm{A}^{}(𝒜)`$, the map $`J^{}:H^{}(X(\widehat{𝒜}))H^{}(X(𝒜))`$, induced by the canonical inclusion, $`J:X(𝒜)X(\widehat{𝒜})`$, corresponds to the canonical projection $`\pi _𝒜^{}:\overline{\mathrm{A}}^{}(𝒜)\mathrm{A}^{}(𝒜)`$.
5. $`\overline{P}_𝒜(T)=P_{\widehat{𝒜}}(T)=_{i=1}^{\mathrm{}}(1+d_iT)`$.
6. There exist minimal CW-complexes $`X`$ and $`Y`$, and homotopy equivalences $`\varphi :XX(𝒜)`$ and $`\psi :X(\widehat{𝒜})Y`$, such that the composite $`\psi J\varphi `$ is homotopic to a cellular classifying map,
$$j:XY,$$
which restricts to the identity on $`p`$-skeleta.
7. The dual of the split exact sequence
$$0H_{}(X)\stackrel{j_{}}{}H_{}(Y)\stackrel{\mathrm{\Pi }_{}}{}H_{}(Y,X)0$$
from (2.2) may be identified with
(4.4)
$$0\mathrm{ker}(\pi _𝒜^{})\stackrel{\iota ^{}}{}\overline{\mathrm{A}}^{}(𝒜)\stackrel{\pi _𝒜^{}}{}\mathrm{A}^{}(𝒜)0,$$
where $`\iota ^{}`$ denotes the inclusion of the kernel.
###### Proof.
(1) The supersolvable arrangement $`\widehat{𝒜}`$ is obtained from $`𝒜`$ by the deformation method introduced in , and refined in . Part (1) follows from this deformation method, which proceeds inductively, using the given composition series of $`𝒜`$.
(2) Up to homotopy, we may view each $`\widehat{𝒜}_i`$ as an arrangement in $`^i`$, and replace each map $`p_i`$ by a bundle map, $`q_i`$, with the specified fiber (more precisely, by a linear fibration, admitting a section, see ). Moreover, the defining polynomials for $`\widehat{𝒜}_i`$ may be written inductively as $`f_1(z_1)=z_1`$, and $`f_i(z_1,\mathrm{},z_i)=f_{i1}(z_1,\mathrm{},z_{i1})_{k=1}^{d_i}(z_ig_{i,k}(z_1,\mathrm{},z_{i1}))`$. Clearly, $`f_i/f_{i1}`$ is a completely solvable Weierstrass polynomial over $`X(\widehat{𝒜}_{i1})`$. Thus, by \[5, Thm. 2.3\], the monodromy of the bundle map $`q_i`$ factors through the pure braid group $`P_{d_i}`$, acting on the free group $`F_{d_i}`$ via the Artin representation.
(3) The first assertion follows from . The specified structure of $`\pi `$ is provided by Part (2).
(4) Since $`\widehat{𝒜}`$ is supersolvable, $`\mathrm{A}^{}(\widehat{𝒜})\overline{\mathrm{A}}^{}(\widehat{𝒜})`$, see Falk and Shelton and Yuzvinsky . Moreover, Theorem 2.4 from insures that $`𝒜`$ and $`\widehat{𝒜}`$ have the same collinearity relations, which implies that $`\overline{\mathrm{A}}^{}(𝒜)\overline{\mathrm{A}}^{}(\widehat{𝒜})`$. The canonical isomorphism in (4) is then given by:
(4.5)
$$H^{}(X(\widehat{𝒜}))\mathrm{A}^{}(\widehat{𝒜})\overline{\mathrm{A}}^{}(\widehat{𝒜})\overline{\mathrm{A}}^{}(𝒜).$$
The identification of $`J^{}`$ with $`\pi _𝒜^{}`$ follows from the fact that the basis $`\{e_1,\mathrm{},e_n\}`$ of $`\mathrm{A}^1`$ is dual to the basis of $`H_1`$ given by the meridians, and $`J_1:H_1(X(𝒜))H_1(X(\widehat{𝒜}))`$ preserves those meridians.
(5) The equality between the Poincaré polynomials of $`\overline{\mathrm{A}}^{}(𝒜)`$ and $`\mathrm{A}^{}(\widehat{𝒜})`$ follows from (4). The second equality follows from , via Part (1).
(6) This follows from Theorem 4.2.
(7) This follows from Parts (4) and (6). ∎
We record as a corollary the most important (for our purposes) consequence of the above theorem.
###### Corollary 4.9.
All hypersolvable complements, and their $`K(\pi ,1)`$ spaces, are minimal, with cohomology algebra generated in degree one.
The following corollary shows that the order of $`\pi _1`$-connectivity of the complement of a hypersolvable arrangement is combinatorially determined (though $`\pi _1`$ itself is not a priori combinatorial).
###### Corollary 4.10.
Let $`𝒜`$ be a hypersolvable arrangement. Set $`X=X(𝒜)`$, and $`\pi =\pi _1(X)`$. Let $`p=p(X)`$ be the order of $`\pi _1`$-connectivity of $`X`$. Then:
1. $`p(X)=sup\{kP_𝒜(T)\overline{P}_𝒜(T)mod(T^{k+1})\}`$.
2. $`p(X)2`$.
3. $`p(X)=\mathrm{}P_𝒜(T)=\overline{P}_𝒜(T)𝒜\text{ is supersolvable}`$.
4. If $`p(X)<\mathrm{}`$, then $`\overline{P}_𝒜(T)P_𝒜(T)c_{p+1}T^{p+1}mod(T^{p+2})`$, where $`c_{p+1}`$ is a positive integer.
###### Proof.
(1) Follows from Theorem 4.8, Parts (3) and (4).
(2) Follows from Remark 2.8 and Theorem 4.8, Part (6).
(3) Follows from (1) and \[22, Prop. 3.4\].
(4) Follows from (1) and the fact that $`\overline{P}_𝒜(T)P_𝒜(T)`$. ∎
### 4.11. A presentation for $`\pi _p(X(𝒜))`$
We come now to the main result in this section. Together with Corollary 4.10(3)–(4), this result provides a complete proof of Theorem 1.4 from the Introduction.
Let $`𝒜`$ be a hypersolvable arrangement. Set $`\pi =\pi _1(X(𝒜))`$ and $`p=p(X(𝒜))`$. Denote by $`\widehat{𝒜}`$ the supersolvable deformation of the associated essential hypersolvable arrangement, $`𝒜_{\mathrm{ess}}`$. Use Theorem 4.8(6) to replace, up to homotopy, the inclusion $`J:X(𝒜_{\mathrm{ess}})X(\widehat{𝒜})`$ by a map, $`j:XY`$, satisfying the hypotheses of Theorem 2.10.
###### Theorem 4.12.
Let $`𝒜`$ be a hypersolvable arrangement, with fundamental group $`\pi =\pi _1(X(𝒜))`$ and order of $`\pi _1`$-connectivity $`p=p(X(𝒜))`$. Then:
1. $`X(𝒜)`$ is aspherical $`p=\mathrm{}`$.
2. If $`p<\mathrm{}`$, then the first non-vanishing higher homotopy group of $`X(𝒜)`$ is $`\pi _p(X(𝒜))`$; as a $`\pi `$-module, $`\pi _p(X(𝒜))`$ is isomorphic to $`\mathrm{coker}(D_p)`$, where the presentation matrix $`D_p`$ is given by (2.3).
3. If $`p<\mathrm{}`$, then the group of $`\pi `$-coinvariants of $`\pi _p(X(𝒜))`$ is free abelian, of rank
$$c_{p+1}=\text{coefficient of }T^{p+1}\text{ in }\overline{P}_𝒜(T)P_𝒜(T).$$
In particular, both $`p`$ and the group $`(\pi _p(X(𝒜)))_\pi `$ are combinatorially determined.
###### Proof.
Parts (1) and (2) follow from Theorem 2.10 and Corollary 2.11. Part (3) follows from Corollary 2.11, Theorem 4.8(7) and Corollary 4.10(1). ∎
###### Remark 4.13.
$`(i)`$ Let $`(C_{}(\stackrel{~}{Y}),_{})`$ be the $`\pi `$-equivariant chain complex of the universal cover, associated to the minimal Morse-theoretic cell structure, $`Y`$, of $`X(\widehat{𝒜})`$, as in Theorem 4.12. Since $`X(\widehat{𝒜})`$ is aspherical, $`(C_{}(\stackrel{~}{Y}),_{})\stackrel{\mathit{ϵ}}{}`$ is a finite, free, $`ϵ`$-minimal, $`\pi `$-resolution of $``$. In particular, it is chain homotopy equivalent to the $`ϵ`$-minimal, finite, free, $`\pi `$-resolution of $``$ constructed by Fox calculus in , starting from the iterated semidirect product structure of $`\pi `$, with trivial monodromy actions on homology, described in Theorem 4.8, Parts (2) and (3).
$`(ii)`$ More generally, let $`j:XY`$ be a classifying map which satisfies the requirements from Theorem 2.10. We can show that the second nilpotent quotient of the first higher non-trivial homotopy group of $`X`$, $`\pi _p(X)/I^2\pi _p(X)`$, is determined, as a filtered module over $`\pi _1(X)/I^2`$, by the map induced by $`j`$ between cohomology rings. For arbitrary hypersolvable complements, the above nilpotent quotient turns out, in this way, to be combinatorially determined. Details will appear elsewhere.
### 4.14. Comparison with some results of Randell
A formula for the coinvariants of the first non-vanishing higher homotopy group, similar to our 4.12(3), was obtained by Randell, using different methods, in \[30, Thm. 2 and Prop. 9\], for the class of iterated generic hyperplane sections (of rank $`3`$) of essential, aspherical arrangements. For an arrangement $`𝒜`$ in this class, $`p(X(𝒜))=\mathrm{rank}(𝒜)1`$, by results from . Randell’s class of arrangements, and the class of hypersolvable arrangements have a similar behavior, from the point of view of the coinvariants of the first higher non-vanishing homotopy group. Nevertheless, the two classes are distinct, as the next two examples show:
###### Example 4.15.
For $`\mathrm{}5`$, let $`𝒜_{\mathrm{}}:=_{\mathrm{}}\{H\}`$, where $`_{\mathrm{}}=\{z_iz_j=0\}_{1i<j\mathrm{}}`$ and $`H=\{z_1+z_2+z_33z_{\mathrm{}}=0\}`$. Each arrangement $`𝒜_{\mathrm{}}`$ is hypersolvable, of rank $`\mathrm{}1`$ and length $`\mathrm{}`$. We claim that $`p(X(𝒜_{\mathrm{}}))=2`$. It follows that these arrangements cannot be iterated generic sections of essential, aspherical arrangements, since this would imply that $`p(X(𝒜_{\mathrm{}}))=\mathrm{}2`$.
The claim may be verified by showing that $`\mathrm{rank}A^3(𝒜_{\mathrm{}})<\mathrm{rank}\overline{A}^3(𝒜_{\mathrm{}})`$; see Corollary 4.10, Parts (1) and (2). Let $`𝒞=\{H_1,H_2,H_3,H\}`$, with $`H_i=\{z_iz_{\mathrm{}}=0\}`$, and let $`\{e_1,e_2,e_3,e\}`$ be the corresponding OS-generators. It is easy to check that $`\mathrm{rank}A^3(𝒜_{\mathrm{}})\mathrm{rank}(\overline{A}(_{\mathrm{}})^{}(e)/(e_𝒞))^3`$, and $`\mathrm{rank}\overline{A}^3(𝒜_{\mathrm{}})=\mathrm{rank}(\overline{A}(_{\mathrm{}})^{}(e))^3`$, directly from the definitions (see §4.3). Notice that $`\{H_1,H_2,H_3\}_{\mathrm{}}`$ is a boolean subarrangement, hence $`e_1e_2e_3`$ is a non-zero element of $`\overline{A}(_{\mathrm{}})`$ (use \[26, Prop. 3.66\]). We infer that $`e_𝒞`$ is a non-zero element of $`\overline{A}(_{\mathrm{}})^{}(e)`$, whence the desired inequality.
###### Example 4.16.
Let $`𝒜`$ be an iterated generic section of an essential, aspherical arrangement $``$ which is not hypersolvable. For example, take $``$ to be the reflection arrangement of type $`\mathrm{D}_n`$, with $`n4`$, see . If $`\mathrm{rank}(𝒜)3`$, then necessarily $`𝒜`$ and $``$ have the same collinearity relations, and therefore $`𝒜`$ cannot be hypersolvable.
There is, however, a certain overlap between the two classes. For instance, iterated generic sections (of rank $`3`$) of fiber-type arrangements are obviously hypersolvable.
###### Example 4.17.
Let $`𝒜`$ be an essential, proper, $`k`$-generic section of an essential arrangement $`\widehat{𝒜}`$, with $`k=\mathrm{rank}(𝒜)1`$. It follows that $`𝒜`$ must be an iterated generic hyperplane section of $`\widehat{𝒜}`$; see the proof of Proposition $`14`$ from . As a particular case, consider a hypersolvable arrangement $`𝒜`$ in $`^3`$, such that $`p(X(𝒜))<\mathrm{}`$ (i.e., such that $`𝒜`$ is not supersolvable, see Corollary 4.10(3)). We must have $`\mathrm{rank}(𝒜)=3`$, since all arrangements are fiber-type in rank $`2`$. Consequently, $`𝒜`$ is a $`2`$-generic, proper section of its essential supersolvable deformation, $`\widehat{𝒜}`$, , hence also an iterated generic section of the aspherical arrangement $`\widehat{𝒜}`$.
## 5. On the structure of $`\pi _2`$ as a $`\pi _1`$-module
We now analyze in more detail the structure of $`\pi _2(X^{})`$, viewed as a module over $`\pi _1(X^{})`$, in the case when $`X^{}=X(𝒜^{})`$ is the complement of an affine line arrangement whose cone $`𝒜=𝐜𝒜^{}`$ is hypersolvable.
### 5.1. $`K(\pi ,1)`$ tests
Let $`𝒜^{}`$ be an arrangement of affine lines in $`^2`$. The complement $`X^{}=X(𝒜^{})`$ has the homotopy-type of a $`2`$-complex, hence the only obstruction to $`X^{}`$ being aspherical is the second homotopy group, $`\pi _2(X^{})`$. In , Falk gave several conditions (some sufficient, some necessary), for the vanishing of $`\pi _2(X^{})`$, providing a (partial) $`K(\pi ,1)`$-test for complexified line arrangements. This test is geometric in nature, involving Gersten-Stallings weight systems.
Another partial $`K(\pi ,1)`$-test, valid this time in all dimensions, but only for hypersolvable arrangements, was given in \[21, Thm. D\]. This test is purely combinatorial. Assuming the cone $`𝒜=𝐜𝒜^{}`$ is hypersolvable, it says that $`X^{}`$ is aspherical if and only if $`\mathrm{}(𝒜)=\mathrm{rank}(𝒜)`$.
None of these asphericity tests, though, gives a precise description of $`\pi _2(X^{})`$, viewed as a $`\pi _1(X^{})`$-module. Our machinery affords such a description, at least in the special case when $`𝒜`$ is hypersolvable.
### 5.2. Affine arrangements with hypersolvable cones
We first describe the structure of the fundamental group of the complement of a deconed fiber-type arrangement, of arbitrary rank.
###### Lemma 5.3.
Let $`𝒜`$ be a supersolvable arrangement, with composition series $`𝒜_1\mathrm{}𝒜_{\mathrm{}}`$, and let $`F_d_{\mathrm{}}_{\rho _\mathrm{}1}F_{d_\mathrm{}1}\mathrm{}_{\rho _2}F_{d_2}_{\rho _1}F_1`$ be the corresponding iterated semidirect product decomposition of $`\pi _1(X(𝒜))`$. If $`𝐝𝒜`$ is a decone of $`𝒜`$, then $`\pi _1(X(𝐝𝒜))=F_d_{\mathrm{}}_{\rho _\mathrm{}1}\mathrm{}_{\rho _2}F_{d_2}`$.
###### Proof.
Recall from the proof of Theorem 4.8(2) that $`𝒜`$ has defining polynomial of the form $`f_𝒜=f_1f_2\mathrm{}f_{\mathrm{}}`$, where $`f_1(z)=z_1`$, and $`f_i/f_{i1}`$ is a completely solvable Weierstrass polynomial over $`X(𝒜_{i1})`$. The decone $`𝐝𝒜`$, obtained by setting $`z_1=1`$, has defining polynomial $`f_{𝐝𝒜}(z_2,\mathrm{},z_{\mathrm{}})=f_2(1,z_2)\mathrm{}f_{\mathrm{}}(1,z_2,\mathrm{},z_{\mathrm{}})`$. The result follows at once. ∎
The above result gives the structure of fundamental groups of complements of deconed hypersolvable arrangements in all dimensions, via the deformation method from Theorem 4.8, Parts (1)–(3) and (6). In dimension $`3`$, this can be much improved, to a precise description of the homotopy type and of the second homotopy group, as follows.
###### Theorem 5.4.
Let $`𝒜^{}`$ be an affine line arrangement in $`^2`$, such that $`𝒜=𝐜𝒜^{}`$ is hypersolvable. Set $`\mathrm{}=\mathrm{}(𝒜)`$, $`X^{}=X(𝒜^{})`$, $`\pi ^{}=\pi _1(X^{})`$, and $`p=p(X^{})`$. Denote by $`𝒜_{\mathrm{ess}}`$ the associated essential hypersolvable arrangement, with supersolvable deformation $`\widehat{𝒜}`$, as in Theorem 4.8. Then $`X^{}`$ has the homotopy type of the $`p`$-skeleton, $`Y^{(p)}`$, of a minimal cell structure, $`Y^{}`$, for $`X(𝐝\widehat{𝒜})`$. In particular:
1. $`X^{}`$ is aspherical $`\mathrm{}3p2.`$
2. If $`\mathrm{}>3`$, then $`\pi _2(X^{})`$ is non-trivial, and admits the following finite, free, $`ϵ`$-minimal $`\pi ^{}`$-resolution:
(5.1)
$$0C_\mathrm{}1\stackrel{_\mathrm{}1}{}C_\mathrm{}2\mathrm{}C_4\stackrel{_4}{}C_3\pi _2(X^{})0,$$
where $`(C_{},_{})`$ is the $`\pi ^{}`$-equivariant chain complex of $`\stackrel{~}{Y^{}}`$.
###### Proof.
If $`p=\mathrm{}`$, then $`\widehat{𝒜}=𝒜_{\mathrm{ess}}`$, by Corollary 4.10(3). Up to homotopy, $`X^{}X(𝐝\widehat{𝒜})`$, and we may use the Morse-theoretic minimal structure from \[9, Corollary 6\]. If $`p<\mathrm{}`$, we know from Example 4.17 that $`𝒜=𝒜_{\mathrm{ess}}`$ is an iterated generic hyperplane section of $`\widehat{𝒜}`$, with $`p=2`$; see also the first paragraph of §4.14. The method of proof of Corollary $`6`$ from provides the desired minimal complex $`Y^{}`$.
(1) The space $`X^{}`$ is aspherical if and only if $`\mathrm{}3`$, by \[21, Thm. D\] (since $`𝒜`$ is fiber-type, if $`\mathrm{}3`$). If $`p2`$, then necessarily $`p>2`$ (by Remark 2.8). It follows from Theorem 4.12(2) that $`\pi _2(X^{})=\pi _2(X(𝒜))=0`$, and so $`X^{}`$ must be aspherical. Conversely, if $`p=2`$, then $`\pi _2(X^{})`$ must be non-zero (use Corollary 2.11).
(2) If $`\mathrm{}>3`$, we know from Part (1) that $`X^{}Y^{(2)}`$, where $`dimY^{}=\mathrm{}1`$. Everything then follows from Remarks 2.12 and 2.13. ∎
###### Remark 5.5.
The resolution (5.1) may have arbitrary length. Indeed, for each $`\mathrm{}1`$, there exists a hypersolvable arrangement $`𝒜`$ in $`^3`$ with $`\mathrm{}(𝒜)=\mathrm{}`$, see \[21, §1\]. Moreover, any sequence of exponents, $`\{1=d_1,d_2,\mathrm{},d_{\mathrm{}}\}`$, may be obtained in this way.
### 5.6. Structure of $`\pi _2`$ of a hypersolvable line arrangement complement
The group of $`\pi ^{}`$-coinvariants of $`\pi _2(X^{})`$ is very simple to describe: By Theorem 5.4(2), it is free abelian, of rank $`b_3(X(𝐝\widehat{𝒜}))=b_3(\pi ^{})`$. On the other hand, the following result shows that $`\pi _2(X^{})`$, when non-trivial, has a fairly complicated structure as a $`\pi ^{}`$-module.
###### Theorem 5.7.
Let $`𝒜^{}`$ be an affine line arrangement in $`^2`$ such that $`𝒜=𝐜𝒜^{}`$ is hypersolvable. Let $`\mathrm{}`$ be the length of $`𝒜`$, and $`\{1=d_1,d_2,\mathrm{},d_{\mathrm{}}\}`$ the exponents. Set $`X^{}=X(𝒜^{})`$ and $`\pi ^{}=\pi _1(X^{})`$. Assume $`\mathrm{}>3`$ (so that $`\pi _2(X^{})0`$). Then:
1. $`\pi _2(X^{})`$ is a projective $`\pi ^{}`$-module if and only if $`\mathrm{}=4`$. In that case, $`\pi _2(X^{})`$ is free, with rank equal to $`b_3(\pi ^{})=d_2d_3d_4`$.
2. $`\pi _2(X^{})`$ is neither finitely generated as an abelian group, nor nilpotent as a $`\pi ^{}`$-module.
###### Proof.
(1) From resolution (5.1), we see that $`\pi _2(X^{})`$ is isomorphic to $`\mathrm{coker}(_4)=\mathrm{im}(_3)C_2`$. If $`\mathrm{}=4`$, then $`\pi _2(X^{})=C_3`$ is a free $`\pi ^{}`$-module, with rank $`b_3(\pi ^{})`$ given by Theorem 4.8(5). If $`\mathrm{}>4`$, then $`\pi _2(X^{})`$ is not projective, by the minimality of (5.1).
(2) Note first that the $`I`$-adic filtration of the group algebra $`\pi ^{}`$ is Hausdorff, in the sense that $`_{k0}I^k=0`$, where $`I=\mathrm{ker}(ϵ:\pi ^{})`$ is the augmentation ideal. This follows from the fact that $`\pi ^{}`$ is an iterated semidirect product of free groups, where all homology monodromy actions are trivial (cf. Lemma 5.3 and Theorem 4.8, Parts (2) and (3)); therefore, $`\pi ^{}`$ is residually torsion-free nilpotent (see ), and so the $`I`$-adic filtration of $`\pi ^{}`$ must be Hausdorff (see ). It follows that the $`I`$-adic filtration of the free $`\pi ^{}`$-module $`C_2`$ is also Hausdorff.
Assume now that either $`\pi _2(X^{})`$ is finitely generated as an abelian group, or nilpotent as a $`\pi ^{}`$-module. It follows that $`I^k\pi _2(X^{})=0`$, for some $`k0`$, and thus $`\pi _2(X^{})`$ must be a nilpotent, non-trivial $`\pi ^{}`$-module. This implies that $`(g_11)\mathrm{}(g_k1)b=0`$, for some $`g_1,\mathrm{},g_k\pi ^{}\{1\}`$, and $`b\pi ^{}\{0\}`$. On the other hand, $`\pi ^{}`$ has no zero-divisors, since $`\pi ^{}`$ is residually torsion-free nilpotent (see ). This gives the desired contradiction, proving (2). ∎
###### Example 5.8.
Let $`𝒜^{}`$ be the affine line arrangement from Figure 2, with defining polynomial $`f_𝒜^{}=z_1z_2(z_11)(z_21)(z_2z_1)`$. Then $`𝒜=𝐜𝒜^{}`$ is an essential $`3`$-slice of the braid arrangement $`_4`$. Hence, $`𝒜`$ is supersolvable, with length $`\mathrm{}=3`$, and exponents $`\{1,2,3\}`$. We then have $`\pi ^{}=F_3F_2`$, and $`X^{}=K(\pi ^{},1)`$.
###### Example 5.9.
Let $`𝒜^{}`$ be the arrangement from Figure 2, with defining polynomial $`f_𝒜^{}=(z_11)(z_1+1)(2z_12z_21)(2z_12z_2+1)(3z_16z_21)(3z_16z_2+1)`$. Then $`𝒜=𝐜𝒜^{}`$ is the arrangement from Fan \[16, §3.I\]. It is readily seen that $`𝒜`$ is hypersolvable, with $`\mathrm{}=4`$, and $`\mathrm{exp}(𝒜)=\{1,2,2,2\}`$. We then have $`\pi ^{}=F_2\times F_2\times F_2`$, and $`\pi _2(X^{})=(\pi ^{})^8`$. Notice that $`V_1(\pi _2(X^{}))=(^{})^6`$.
###### Example 5.10.
Let $`𝒜^{}`$ be the arrangement from Figure 2, with defining polynomial $`f_𝒜^{}=z_1z_2(z_11)(z_2z_11)(z_2+z_12)`$. Then $`𝒜=𝐜𝒜^{}`$ is hypersolvable, with $`\mathrm{}=5`$, and $`\mathrm{exp}(𝒜)=\{1,1,1,1,2\}`$. The algorithm from yields a “braid monodromy” presentation for $`\pi ^{}`$, with generators $`x_1,\mathrm{},x_5`$, and commutation relations $`[x_i,x_j]=1`$, for $`1i3`$ and $`i<j`$. By , $`X^{}`$ has the homotopy type of the $`2`$-complex associated to this presentation. In turn, this $`2`$-complex is homotopy equivalent to the $`2`$-skeleton of $`Y^{}=T^3\times (S^1S^1)`$, the product of the $`3`$-torus with a wedge of two circles, endowed with the standard minimal cell structure. From Remarks 2.12 and 2.13, we infer that $`\pi _2(X^{})=\mathrm{coker}(_4)`$, where $`(C_{},_{})`$ is the $`\pi ^{}`$-equivariant chain complex of $`\stackrel{~}{Y^{}}`$. The Künneth formula and covering space theory lead to the following presentation of $`\pi _2(X^{})`$, viewed as a left $`\pi ^{}`$-module:
$$(\pi ^{})^2\stackrel{\left(\begin{array}{ccccccc}1x_4& 1x_3& 0& x_21& 0& 1x_1& 0\\ 1x_5& 0& 1x_3& 0& x_21& 0& 1x_1\end{array}\right)}{}(\pi ^{})^7\pi _2(X^{})0.$$
Notice that $`V_6(\pi _2(X^{}))=\{t(^{})^5t_1=t_2=t_3=1\}`$ is a $`2`$-dimensional subtorus. That $`M=\pi _2(X^{})`$ is not nilpotent can be seen directly, as follows. Let $`\stackrel{~}{M}=\pi _2(X^{})_\pi ^{}^5`$, and let $`\mathrm{gr}\stackrel{~}{M}`$ be the associated graded module (with respect to the $`I`$-adic filtration). From the above presentation, we may easily compute its Hilbert series: $`\mathrm{Hilb}(\mathrm{gr}\stackrel{~}{M},t)=\frac{72t}{(1t)^5}`$. Since this series is not a polynomial, $`\stackrel{~}{M}`$ is not nilpotent, and so $`M`$ isn’t, either.
## 6. Graphic arrangements
In this section, we apply our methods to graphic arrangements. We start by giving a graph-theoretic characterization of hypersolvable arrangements within this class. We then show how to get information on $`\pi _2`$ of the complement directly from the graph, in the case of arrangements associated to graphs without $`3`$-cycles.
### 6.1. Graphs and arrangements
Let $`G=(𝒱,)`$ be a non-empty subgraph of the complete graph on a finite set of vertices $`𝒱`$. Assume that there are no isolated vertices in the graph, so that the set of edges $``$ determines $`G`$. All graphs considered in this section will be of this type.
Let $`𝒱=\{1,\mathrm{},m\}`$. The graphic arrangement associated to $`G=(𝒱,)`$ is the arrangement in $`^m`$ given by $`𝒜_G=\{\mathrm{ker}(z_iz_j)\{i,j\}\}`$, see . For each edge $`e=\{i,j\}`$, we will denote by $`H_e:=\mathrm{ker}(z_iz_j)`$ the corresponding hyperplane of $`𝒜_G`$.
Clearly, an arrangement is graphic if and only if it is a sub-arrangement of a braid arrangement. For example, if $`G`$ is the complete graph on $`m`$ vertices, then $`𝒜_G=_m`$, the braid arrangement in $`^m`$. If $`G`$ is a diagram of type $`\mathrm{A}_m`$, then $`𝒜_G`$ is a Boolean arrangement. If $`G`$ is an $`m`$-cycle, then $`𝒜_G`$ is the cone of a generic arrangement.
Many of the usual invariants associated to $`𝒜_G`$ can be computed directly from $`G`$. For example, $`P_{𝒜_G}(T)=(T)^m\chi _G(T^1)`$, where $`\chi _G(T)`$ is the chromatic polynomial of $`G`$, see . Also, an nbc-basis for $`\mathrm{A}^{}(𝒜_G)`$ corresponds to an nbc-basis for $`G`$, as follows. Fix an ordering on the edges, $`_G=\{e_1<\mathrm{}<e_n\}`$, and denote by $`\alpha _i`$ the defining equation of $`H_{e_i}`$. Then, $`\{\alpha _{i_1},\mathrm{},\alpha _{i_r}\}`$ is minimally dependent if and only if $`\{e_{i_1},\mathrm{},e_{i_r}\}`$ is an $`r`$-cycle of $`G`$. Deleting the highest edge from this cycle yields a broken circuit. The resulting nbc-basis for $`\mathrm{A}^{}(𝒜_G)`$ is given by
(6.1)
$$\{e_KK\text{ is a subgraph of }G\text{ which does not contain any broken circuit of }G\},$$
where $`e_K:=e_{i_1}\mathrm{}e_{i_s}^s(e_1,\mathrm{},e_n)`$, if $`_K=\{e_{i_1},\mathrm{},e_{i_s}\}`$.
### 6.2. Supersolvable and hypersolvable graphs
The following results, due to Stanley and Fulkerson and Gross , tell us how to (easily) recognize supersolvable arrangements within the class of graphic arrangements.
###### Theorem 6.3.
Let $`𝒜_G`$ be a graphic arrangement. Then:
(Stanley ) $`𝒜_G`$ is supersolvable if and only if the graph $`G`$ is supersolvable, i.e., it has a supersolvable composition series of induced subgraphs, $`\mathrm{}=G_0G_1\mathrm{}G_{\mathrm{}}=G`$, such that:
1. for each $`1i\mathrm{}`$, there is a single vertex in $`G_iG_{i1}`$, say, $`v_i`$;
2. the subgraph of $`G_i`$ induced by $`v_i`$ and its neighbors in $`G_i`$ is complete.
(Fulkerson and Gross ) $`G`$ has a supersolvable composition series if and only if every cycle in $`G`$ of length greater than $`3`$ has a chord.
As a simple example, consider the two graphic arrangements given by the graphs in Figure 3. Neither graph has a $`3`$-cycle; each graph has $`4`$-cycles, with no chords. Hence, the two arrangements are not supersolvable.
Stanley’s supersolvable test from Theorem 6.3 has a hypersolvable analogue. To state it, we first need a definition.
###### Definition 6.4.
A pair of graphs, $`(G,K)`$, is called a solvable extension if $`K`$ is a subgraph of $`G`$, with $`\mathrm{}_K_G`$, and:
1. There is no $`3`$-cycle in $`G`$ having two edges from $`_K`$ and one edge from $`_G_K`$.
2. Either $`_G_K=\{e\}`$, and both endpoints of $`e`$ are not in $`𝒱_K`$, or there exist distinct vertices, $`\{v_1,\mathrm{},v_k,v\}𝒱_G`$, with $`\{v_1,\mathrm{},v_k\}𝒱_K`$, such that:
1. $`K`$ contains the complete graph on $`\{v_1,\mathrm{},v_k\}`$, and
2. $`_G_K=\{\{v,v_s\}1sk\}`$.
###### Lemma 6.5.
An extension of graphs, $`(G,K)`$, is solvable if and only if the corresponding extension of graphic arrangements, $`(𝒜_G,𝒜_K)`$, is solvable.
###### Proof.
Let $`e_1=\{i_1,j_1\}`$, $`e_2=\{i_2,j_2\}`$, $`e_3=\{i_3,j_3\}`$ be three distinct edges of $`G`$. Notice that the corresponding defining equations, $`\{z_{i_r}z_{j_r}1r3\}`$, viewed as points in $`(^m)`$, are collinear if and only if $`\{e_r1r3\}`$ are the edges of a $`3`$-cycle. Using this remark, it is a straightforward exercise to translate conditions (I)–(III) from §4.4 into conditions (1) and (2) from Definition 6.4. ∎
###### Definition 6.6.
A graph $`G`$ is called hypersolvable if it has a hypersolvable composition series, i.e., a chain of subgraphs, $`G_1\mathrm{}G_iG_{i+1}\mathrm{}G_{\mathrm{}}=G`$, such that $`G_1`$ has a single edge, and $`(G_{i+1},G_i)`$ is a solvable extension, for $`i=1,\mathrm{},\mathrm{}1`$.
The class of hypersolvable graphs contains the supersolvable (or chordal) graphs described in Theorem 6.3, and many others. For example, the graphs from Figure 3 are both hypersolvable, with composition series $`G_i=\{e_1,\mathrm{},e_i\}`$, $`1i9`$, but not supersolvable.
###### Proposition 6.7.
A graph $`G`$ is hypersolvable if and only if the graphic arrangement $`𝒜_G`$ is hypersolvable.
###### Proof.
Clearly, $`G_1\mathrm{}G_{\mathrm{}}`$ is a composition series for $`G`$ if and only if $`𝒜_{G_1}\mathrm{}𝒜_G_{\mathrm{}}`$ is a composition series for $`𝒜_G`$. ∎
### 6.8. Graphs with no $`3`$-cycles
We now analyze in more detail a very simple example: hypersolvable arrangements coming from graphs without $`3`$-cycles, and their second homotopy group.
###### Proposition 6.9.
Let $`G`$ be a graph with no $`3`$-cycles, and with edges $`\{e_1,\mathrm{},e_n\}`$. Then:
1. The graph $`G`$ is hypersolvable, with composition series $`G_i=\{e_1,\mathrm{},e_i\}`$, $`1in`$.
2. The arrangement $`𝒜=𝒜_G`$ is hypersolvable, with length $`n`$ and exponents $`\{1,\mathrm{},1\}`$.
3. $`\overline{\mathrm{A}}^{}(𝒜)=^{}(e_1,\mathrm{},e_n)`$.
4. $`\pi _1(X(𝒜))=^n`$.
###### Proof.
There are no collinearity relations among the defining equations of $`𝒜`$, since $`G`$ has no $`3`$-cycles. Part (1) then follows from Definitions 6.4 and 6.6, Part (2) from Proposition 6.7 and (1), and Part (3) from the definition of the quadratic OS-algebra.
Now set $`m=\mathrm{\#}\{\text{vertices of }G\}`$. If $`m3`$, then Part (4) is trivially verified. If $`m>3`$, we may take a generic $`3`$-plane $`P`$ in $`^m`$ with the property that $`\pi _1(X(𝒜))=\pi _1(X(𝒜P))`$ (by ), and such that $`𝒜`$ and $`𝒜P`$ have the same collinearity relations. The decone of the arrangement $`𝒜P`$ is thus generic, and so $`\pi _1(X(𝒜P))=^n`$ (by ). ∎
###### Theorem 6.10.
Let $`G`$ be a graph with no $`3`$-cycles. Let $`𝒮`$ be the set of $`4`$-cycles of $`G`$. Set $`𝒜=𝒜_G`$ and $`X=X(𝒜)`$. Then:
* If $`𝒮=\mathrm{}`$, then $`\pi _2(X)=0`$.
* If $`𝒮\mathrm{}`$, then $`\pi _2(X)`$ is non-trivial, with group of coinvariants isomorphic to $`[𝒮]`$, the free abelian group generated by $`𝒮`$.
###### Proof.
Set $`n=|𝒜|`$, $`\pi =\pi _1(X)`$ and $`p=p(X)`$. Recall from §4.3 the construction of the OS-algebra of $`𝒜`$, together with the graphic counterpart from §6.1. From Proposition 6.9, we know that $`𝒜`$ hypersolvable, $`\overline{\mathrm{A}}^{}(𝒜)=^{}(e_1,\mathrm{},e_n)`$, and $`\pi =^n`$.
If $`𝒮=\mathrm{}`$, then $`p3`$, by Corollary 4.10(1), and thus $`\pi _2(X)=0`$, by Theorem 4.12(1)-(2).
If $`𝒮\mathrm{}`$, the same argument shows that $`p=2`$. From Theorem 4.12(3) and (4.4), we deduce that $`\pi _2(X)_\pi `$ is isomorphic to $`\mathrm{ker}(\pi _𝒜^3)`$, which is generated by $`\{e_SS𝒮\}`$. A quick inspection of the construction of the nbc-basis in degree $`3`$ reveals that $`\mathrm{ker}(\pi _𝒜^3)`$ is free abelian, with rank equal to the number of broken $`4`$-circuits. Now, in the case of graphic arrangements, the associated broken circuit uniquely determines a cycle in the graph. Hence, the map $`:[𝒮]\mathrm{ker}(\pi _𝒜^3)`$ is actually an isomorphism. ∎
###### Example 6.11.
Consider the hypersolvable graphs without $`3`$-cycles, $`G_1`$ and $`G_2`$, in Figure 3, and let $`𝒜_1`$ and $`𝒜_2`$ be the corresponding graphic arrangements in $`^7`$, with complements $`X_1`$ and $`X_2`$. Both complements have $`\pi _1=^9`$ and $`b_2=\left(\genfrac{}{}{0pt}{}{9}{2}\right)=36`$. The $`4`$-cycles of $`G_1`$ are $`\{e_1,e_2,e_9,e_7\}`$ and $`\{e_5,e_6,e_7,e_8\}`$, while those of $`G_2`$ are $`\{e_1,e_2,e_9,e_7\}`$, $`\{e_1,e_8,e_6,e_7\}`$ and $`\{e_2,e_9,e_6,e_8\}`$. It follows that the spaces $`X_1`$ and $`X_2`$ have distinct homotopy $`2`$-types, since the ranks of the coinvariants of the $`^9`$-modules $`\pi _2(X_1)`$ and $`\pi _2(X_2)`$ are different. |
warning/0002/math0002166.html | ar5iv | text | # 1 INTRODUCTION
## 1 INTRODUCTION
This paper is in some respects a sequel to , which considered group doublecross products, i.e. a group factoring into two subgroups. Group doublecross products are the foundation of one way to look at certain integrable field theories . Here the space-time is imbedded in the group by a function called the ‘classical vacuum map’. This imbedding possibly encodes information about the geometry of the space-time, and in general there is no reason why its image should be a subgroup (also see ). This raised the possibility of considering more general factorisations of groups, and their corresponding algebras. These algebras turned out to be non-trivially associated, after the manner of . In retrospect it is suprising how many of the results of the standard theory of doublecross products and bicrossproducts carry over into the present setting. Typically all that is required is the insertion of a few additional pieces into the relevant formulae. Throught the paper I have liberally used as a reference for tensor categories and braided Hopf algebras. The reader should note that the coset representative constructions in this paper are essentially those appearing in the group cohomology analysis of exact sequences of groups . This then leads on to non-commutative topological cohomology via crossed modules . It is not clear to me whether there is a more direct link between non-trivially associated tensor categories and cohomology.
Take a representative element of every coset for the left action of a subgroup $`G`$ on a group $`X`$, and form a set $`M`$ of these elements. From the algebraic structure on $`X`$ we can construct a binary operation on $`M`$ which has a left identity and right division. Conversely any set with such a binary operation can be realised as a set of left coset representatives for the quotient of two permutation groups. This binary operation is not associative, but the breakdown of associativity is given by a ‘cocycle’ $`\tau :M\times MG`$. Using this cocycle we can construct a non-trivial associator for a category $`𝒞`$ of $`M`$-graded right representations of $`G`$. This category also has evaluation and coevaluation maps, making it into a rigid tensor category. If we assume that the binary operation on $`M`$ has left division (which is not always true) then the grading and group action can be combined into the action of an algebra $`H`$ on the objects in the category. It turns out that $`H`$ itself is in $`𝒞`$, and that the multiplication is associative (using the associator).
Next we construct a double factorisation from two copies of the original group $`X`$, though we give one copy of $`X`$ a different binary operation. The category $`𝒟`$ constructed from the double (in the same way that $`𝒞`$ was constructed from $`X`$) is braided, as well as being non-trivially associated. We can again form an algebra $`D`$ in $`𝒟`$ whose action combines both the gradings and actions in the definition of $`𝒟`$. But now, using the braiding in $`𝒟`$, we can reconstruct a coproduct on $`D`$ from the tensor product structure in $`𝒟`$. The existence of such a braided Hopf algebra is guaranteed by a general reconsruction theory , and we explicitly calculate the braided Hopf algebra structure on $`D`$ using these methods. I do not know whether the braided category $`𝒟`$ gives any interesting knot invariants.
Except for the section on left coset representatives, I shall assume that all groups are finite and that all vector spaces are finite dimensional. This is to avoid problems with measurability and continuity.
I would like to thank Y. Bespalov, Ronnie Brown and Shahn Majid for their assistance in the preparation of this paper.
## 2 Left coset representatives.
###### Definition 2.1
Given a group $`X`$ and a subgroup $`G`$, call $`MX`$ a set of left coset representatives if for every $`xX`$ there is a unique $`sM`$ so that $`xGs`$. We shall call the decomposition $`x=us`$ for $`uG`$ and $`sM`$ the unique factorisation of $`x`$.
For the remainder of this section we assume that $`MX`$ is a set of left coset representatives for the subgroup $`GX`$. The identity in $`X`$ will be denoted $`e`$.
###### Definition 2.2
Given $`s,tM`$, we define $`\tau (s,t)G`$ and $`stM`$ by the unique factorisation $`st=\tau (s,t)(st)`$ in $`X`$. We also define functions $`:M\times GG`$ and $`:M\times GM`$ by the unique factorisation $`su=(su)(su)`$ for $`s,suM`$ and $`u,suG`$.
###### Proposition 2.3
The binary operation $`(M,)`$ has a unique left identity $`e_mM`$ (i.e. $`e_mt=t`$ for all $`tM`$) and has the right division property (i.e. for all $`t,sM`$ there is a unique solution $`pM`$ to the equation $`ps=t`$). If $`eM`$ then $`e_m=e`$ is also a right identity.
Proof There is a unique factorisation $`e=ue_m`$ for $`e_mM`$ and $`uG`$, so $`GM=\{e_m\}`$. Then $`e_mt=t`$ by definition. Conversely if $`st=t`$ then $`st=\tau (s,t)t`$, so $`sGM`$.
If $`ps=t`$ then $`ps=\tau (p,s)t`$, so $`\tau (p,s)^1p=ts^1`$. Now apply unique factorisation to $`ts^1X`$. $``$$``$
By applying the right division property to solve the equation $`pt=e_m`$ for a given $`tM`$, we see that there is a unique left inverse $`t^L`$ for every $`t`$, satisfying the equation $`t^Lt=e_m`$.
We shall use the result of the next proposition at many places in the paper:
###### Proposition 2.4
The following identities between $`(M,)`$ and $`\tau `$ hold, where we take $`t,s,pM`$ and $`u,vG`$:
$`s(tu)=\tau (s,t)((st)u)\tau (s(tu),tu)^1`$ $`\mathrm{and}`$ $`(st)u=(s(tu))(tu),`$
$`suv=(su)((su)v)`$ $`\mathrm{and}`$ $`suv=(su)v,`$
$`\tau (p,s)\tau (ps,t)=(p\tau (s,t))\tau (p\tau (s,t),st)`$ $`\mathrm{and}`$ $`(p\tau (s,t))(st)=(ps)t.`$
Proof We can deduce these identities from the associativity of $`X`$. From the identity $`(st)u=s(tu)`$ we can deduce that
$`(st)u`$ $`=`$ $`\tau (s,t)(st)u=\tau (s,t)((st)u)((st)u),`$
$`s(tu)`$ $`=`$ $`s(tu)(tu)=(s(tu))(s(tu))(tu)`$
$`=`$ $`(s(tu))\tau (s(tu),tu)((s(tu))(tu)).`$
The first line follows by uniqueness of factorisation. From $`s(uv)=(su)v`$,
$`s(uv)`$ $`=`$ $`(suv)(suv),`$
$`(su)v`$ $`=`$ $`(su)(su)v=(su)((su)v)((su)v),`$
giving the second line identities. Finally from $`p(st)=(ps)t`$,
$`p(st)`$ $`=`$ $`p\tau (s,t)(st)=(p\tau (s,t))(p\tau (s,t))(st)`$
$`=`$ $`(p\tau (s,t))\tau (p\tau (s,t),st)((p\tau (s,t))(st)),`$
$`(ps)t`$ $`=`$ $`\tau (p,s)(ps)t=\tau (p,s)\tau (ps,t)((ps)t),`$
giving the last line. $``$$``$
###### Proposition 2.5
The following identities between $`(M,)`$ and $`\tau `$ hold, for all $`tM`$ and $`vG`$:
$`e_mv=e_m,e_mv=e_mve_m^1`$ , $`te=e,te=t,`$
$`\tau (e_m,t)=e_m,te_m^1`$ $`=`$ $`\tau (te_m^1,e_m)^1,(te_m^1)e_m=t.`$
Proof We have the factorisation $`e_mv=(e_mve_m^1)e_m`$, where $`e_mM`$ and $`e_mve_m^1G`$. Also $`te=et`$ for $`eG`$ and $`tM`$. Next $`e_mt=\tau (e_m,t)(e_mt)=\tau (e_m,t)t`$. Finally
$$t=te_m^1e_m=(te_m^1)(te_m^1)e_m=(te_m^1)\tau (te_m^1,e_m)((te_m^1)e_m),$$
giving the last identities. $``$$``$
This last proposition makes sense because $`e_mGM`$. For situations where it is convenient to forget about the original group $`X`$, and just concentrate on $`G`$ and $`(M,)`$, we will use $`e_mM`$ for the left identity in $`(M,)`$, and set $`f_m=e_mG`$.
###### Example 2.6
Take $`X`$ to be the permutation group $`S_3`$ of 3 objects $`\{1,2,3\}`$, and let $`G`$ be the non-normal subgroup $`\{e,(12)\}`$. Take the set of left coset representatives $`M=\{(12),(13),(23)\}`$. The dot and $`\tau `$ operation are given by the following tables, where the row $`s`$ column $`t`$ entry corresponds to $`st`$ or $`\tau (s,t)`$:
$``$ (12) (13) (23) (12) (12) (13) (23) (13) (23) (12) (13) (23) (13) (23) (12) $`\tau `$ (12) (13) (23) (12) (12) (12) (12) (13) (12) (12) (12) (23) (12) (12) (12)
The fact that $`(M,)`$ satisfies the right division property is just the condition that every element of $`M`$ appears exactly once in each column of the table for dot. In this case we also see that every element of $`M`$ appears exactly once in each row of the $`(M,)`$ table, so $`(M,)`$ satisfies left division. However there is no right identity, so $`(M,)`$ does not form a group.
###### Example 2.7
Take $`X`$ to be the permutation group $`S_3`$ of 3 objects $`\{1,2,3\}`$, and let $`G`$ be the non-normal subgroup $`\{e,(12)\}`$. Take the set of left coset representatives $`M=\{e,(23),(13)\}`$. The operation $``$ is trivial, and $``$ is given by the action of $`(12)`$ on $`M`$ swapping $`(23)`$ and $`(13)`$. The dot and $`\tau `$ operation are given by the following tables:
$``$ e (13) (23) e e (13) (23) (13) (13) e (13) (23) (23) (23) e $`\tau `$ e (13) (23) e e e e (13) e e (12) (23) e (12) e
This time we see that $`(M,)`$ does not satisfy left division, but does have a 2-sided identity.
###### Example 2.8
Take $`X`$ to be the permutation group $`S_3`$ of 3 objects $`\{1,2,3\}`$, and let $`G`$ be the non-normal subgroup $`\{e,(12)\}`$. Take the set of left coset representatives $`M=\{e,(123),(132)\}`$. In this case $`M`$ is a subgroup, and the subgroup operation is the dot product. This is just the case of a group doublecross product .
###### Example 2.9
Take $`X`$ to be the dihedral group $`D_6=x,y:x^6=y^2=e,xy=yx^5`$, and $`G`$ to be the non-abelian normal subgroup of order 6 generated by $`x^2`$ and $`y`$. We choose $`M=\{e,x\}`$. The dot operation on $`M`$ is given by $`e`$ the 2-sided identity and $`xx=e`$. The $`\tau `$ function is given by $`\tau (x,x)=x^2`$, and all other combinations giving $`e`$. The operation $``$ is trivial, and $``$ is given by $`x`$ acting on $`G=\{e,x^2,x^4,y,yx^2,yx^4\}`$ as the permutation $`(y,yx^4,yx^2)`$, i.e. $`xy=yx^4`$ etc. Observe that though $`(M,)`$ is a group, $``$ is not a group action.
###### Example 2.10
Take $`X`$ to be the group $`S_5`$ of permutations of the objects $`\{1,2,3,4,5\}`$, and $`G`$ to be the subgroup fixing the object $`1`$. We choose $`M=\{e,(12)(354),(14253),(15234),(13245)\}`$. If we set $`a=(12)(354)`$, $`b=(14253)`$, $`c=(15234)`$ and $`d=(13245)`$, we get the tables:
$``$ e a b c d e e a b c d a a e c d b b b c d a e c c d e b a d d b a e c $`\tau `$ e a b c d e e e e e e a e (345) (2534) (2345) (2453) b e (34) (354) (2345) (354) c e (45) (354) (254) (2453) d e (35) (2534) (354) (235)
We see that $`(M,)`$ satisfies right and left division, that $`e`$ is a 2-sided identity, but that $`(M,)`$ is not a group.
The last example is just an application of a general construction:
###### Theorem 2.11
: A modified Cayley’s theorem. Any set $`F`$ with a binary operation $``$ which has a left identity and right division can be imbedded in $`S_F`$ (the group of permutations of the elements of $`F`$), as a set of left coset representatives for the subgroup $`GS_F`$ which fixes the left identity.
Proof The function $`\sigma :FS_F`$ is defined by $`\sigma (g)(fg)=f`$. Note that $`\sigma (g)`$ is a 1-1 correspondence because $`(F,)`$ has right division.
Let $`e_F`$ be the left identity in $`(F,)`$, which is unique by right division. Take any $`\varphi S_F`$, and set $`g=\varphi ^1(e_F)`$. Then $`\psi =\varphi \sigma (g)^1`$ has the property that $`\psi (e_F)=e_F`$, i.e. $`\psi G`$. Further if $`\chi \sigma (g)=\sigma (f)`$ for any $`\chi G`$, then by applying $`\sigma (g)^1\chi ^1=\sigma (f)^1`$ to $`e_F`$ we see that $`g=f`$. We conclude that the image of $`\sigma `$ forms a set of left coset representatives for the subgroup $`G`$.
Now consider the equation $`\sigma (f)\sigma (g)=\chi \sigma (h)`$, for $`\chi G`$. Applying the inverse of each side to $`e_F`$ we see $`\sigma (g)^1(\sigma (f)^1(e_F))=\sigma (g)^1(f)=fg`$, so $`h=fg`$ as required. $``$$``$
###### Proposition 2.12
The subgroup $`G`$ is normal in $`X`$ if and only if $``$ is trivial, i.e. $`su=s`$ for all $`sM`$ and $`uG`$. If $`G`$ is normal then $`(M,)`$ is isomorphic to the quotient group $`G\backslash X`$, the isomorphism being the restriction of the quotient map $`XG\backslash X`$. The subset $`M`$ is a subgroup of $`X`$ if and only if $`eM`$ and $`\tau (s,t)=e`$ for all $`s,tM`$.
Proof Since $`G`$ is closed under conjugation by elements of $`G`$, we just have to check conjugation by elements of $`M`$ to see if $`G`$ is normal. Then for all $`sM`$ and $`uG`$:
$$sus^1=(su)(su)s^1GsuGssu=s.$$
If $`G`$ is normal then we just use the usual definition of multiplication of left cosets.
If $`\tau (s,t)=e`$ for all $`s,tM`$, then $`st=\tau (s,t)(st)=stM`$, so the subset $`M`$ is closed under multiplication in $`X`$. If in addition $`eM`$ then for every $`tM`$ (by right division) there is a $`t^LM`$ so that $`t^Lt=e`$. Then $`t^Lt=\tau (t^L,t)(t^Lt)=e`$, so $`t^L=t^1`$, and $`M`$ is closed under inverse in $`X`$. $``$$``$
We would like to remove the dependence on the group $`X`$, and say that certain conditions on $`G`$, $`(M,)`$, $`\tau `$, $``$ and $``$ are equivalent to the existence of the group $`X`$. To this end, for the remainder of this section we forget how $`(M,)`$ was constructed, and just begin with a group $`G`$ and a set with binary operation $`(M,)`$.
###### Proposition 2.13
Suppose that the functions $`:M\times GG`$, $`:M\times GM`$ and $`\tau :M\times MG`$ satisfy the identities in (2.4). Then the binary operation on the set $`G\times M`$ defined by $`(u,s)(v,t)=(u(sv)\tau (sv,t),(sv)t)`$ is associative.
Proof Begin by calculating
$`((u,s)(v,t))(w,p)`$ $`=`$ $`(u(sv)\tau (sv,t)((sv)tw)\tau (((sv)t)w,p),(((sv)t)w)p),`$ (1)
$`(u,s)((v,t)(w,p))`$ $`=`$ $`(u,s)(v(tw)\tau (tw,p),(tw)p)`$ (2)
$`=`$ $`(u(sv(tw)\tau (tw,p))\tau (sv(tw)\tau (tw,p),(tw)p)`$ (4)
$`,(sv(tw)\tau (tw,p))((tw)p)).`$
To show that the $`M`$ components of (4) are identical, we use
$$(((sv)t)w)p=((sv(tw))(tw))p=(sv(tw)\tau (tw,p))((tw)p).$$
We can use the identities to show that the $`G`$ component of $`((u,s)(v,t))(w,p)`$ is
$$u(sv)((sv)(tw))\tau (sv(tw),tw)\tau ((sv(tw))(tw),p),$$
wheras the $`G`$ component of $`(u,s)((v,t)(w,p))`$ is
$$u(sv)((sv)(tw))((sv(tw))\tau (tw,p))\tau (sv(tw)\tau (tw,p),(tw)p),$$
and then use the identities again to show that these are the same. $``$$``$
###### Proposition 2.14
Suppose that the functions $`:M\times GG`$, $`:M\times GM`$ and $`\tau :M\times MG`$ satisfy the identities in (2.4). Suppose that there is a left identity $`e_mM`$ for $`(M,)`$ and an element $`f_mG`$ so that for all $`tM`$ and $`vG`$,
$`e_mv=e_m,e_mv=f_mvf_m^1`$ , $`te=e,te=t,`$
$`\tau (e_m,t)=f_m,tf_m^1`$ $`=`$ $`\tau (tf_m^1,e_m)^1,(tf_m^1)e_m=t.`$
Then the multiplication on the set $`G\times M`$ defined in the previous proposition has a 2-sided identity $`(f_m^1,e_m)`$.
If in addition we suppose that $`(M,)`$ has left inverses (i.e. for every $`tM`$ there is a $`t^LM`$ so that $`t^Lt=e_m`$), then $`G\times M`$ has left inverses, defined by
$$(v,t)^L=(f_m^1\tau (t^L,t)^1(t^Lv^1),t^Lv^1).$$
These properties imply that $`G\times M`$ with the given structure is a group.
Proof First we check the 2-sided identity
$`(f_m^1,e_m)(v,t)`$ $`=`$ $`(f_m^1(e_mv)\tau (e_mv,t),(e_mv)t)`$
$`=`$ $`(f_m^1f_mvf_m^1f_m,e_mt)=(v,t),`$
$`(v,t)(f_m^1,e_m)`$ $`=`$ $`(v(tf_m^1)\tau (tf_m^1,e_m),(tf_m^1)e_m)=(v,t).`$
Finally we check the left inverse:
$`(f_m^1\tau (t^L,t)^1(t^Lv^1),t^Lv^1)(v,t)`$ $`=`$ $`(f_m^1\tau (t^L,t)^1(t^Lv^1)((t^Lv^1)v)\tau (t^L,t),t^Lt)`$
$`=`$ $`(f_m^1\tau (t^L,t)^1(t^Lv^1v)\tau (t^L,t),e_m)=(f_m^1,e_m).`$
It is now standard algebra to check that these conditions on identities and inverses, together with associativity, give a group structure. $``$$``$
We can now imbed the group $`G`$ in $`G\times M`$ by the map $`v(vf_m^1,e_m)`$, and $`M`$ in $`G\times M`$ by the map $`t(e,t)`$. Then we get the original situation with left coset representatives.
## 3 A tensor category
Take a group $`X`$ with subgroup $`G`$, and a set of left coset representatives $`M`$. We again take $`e_m`$ to be the left identity in $`M`$ and $`f_m`$ to be the corresponding element in $`G`$. If the reader wishes, the situation can be simplified if $`eM`$, as then $`e_m=f_m=e`$.
Take a category $`𝒞`$ of finite dimensional vector spaces over a field $`k`$, whose objects are right representations of the group $`G`$ and possess $`M`$-gradings, i.e. an object $`V`$ decomposes as a direct sum of subspaces $`V=_{sM}V_s`$. If $`\xi V_s`$ for some $`sM`$ we say that $`\xi `$ is a homogenous element of $`V`$, with grade $`\xi =s`$. In our formulae in this paper we shall usually assume that we have chosen homogenous elements of the relevant objects, as the general elements are just linear combinations of the homogenous elements. We write the action for the representation as $`\overline{}:V\times GV`$. In addition we suppose that the action and grading satisfy a compatibility condition, $`\xi \overline{}u=\xi u`$. The morphisms are linear maps which preserve both the grading and the action, i.e. for a morphism $`\theta :V\stackrel{~}{V}`$ we have $`\theta (\xi )=\xi `$ and $`\theta (\xi )\overline{}u=\theta (\xi \overline{}u)`$ for all $`\xi V`$ and $`uG`$.
###### Proposition 3.1
We can make $`𝒞`$ into a tensor category by taking $`VW`$ to be the usual vector space tensor product, with actions and gradings given by
$$\xi \eta =\xi \eta \mathrm{and}(\xi \eta )\overline{}u=\xi \overline{}(\eta u)\eta \overline{}u.$$
For morphisms $`\theta :V\stackrel{~}{V}`$ and $`\varphi :W\stackrel{~}{W}`$ we define the morphism $`\theta \varphi :VW\stackrel{~}{V}\stackrel{~}{W}`$ by $`(\theta \varphi )(\xi \eta )=\theta (\xi )\varphi (\eta )`$, which is just the usual vector space formula.
Proof We must check that $`(\xi \eta )\overline{}u=\xi \eta u`$, which is automatic from the usual identities (2.4). Also we have to check that $`((\xi \eta )\overline{}u)\overline{}v=(\xi \eta )\overline{}uv`$, which is again simple from the identities. It is also straightforward to check that $`\theta \varphi `$ is a morphism in $`𝒞`$. $``$$``$
###### Proposition 3.2
The identity for the tensor operation is just the vector space $`k`$ with trivial $`G`$-action and grade $`e_mM`$. For any object $`V`$ the morphisms $`l_V:VVk`$ and $`r_V:VkV`$ are given by the formulae $`l_V(\xi )=\xi \overline{}f_m^11`$ and $`r_V(\xi )=1\xi `$, where $`1`$ is the multiplicative identity in $`k`$.
Proof We must check that the maps $`l_V`$ and $`r_V`$ are morphisms in $`𝒞`$. We have
$`(\xi \overline{}f_m^11)\overline{}u`$ $`=`$ $`\xi \overline{}f_m^1(1u)1\overline{}u=\xi \overline{}f_m^1(e_mu)1=\xi \overline{}uf_m^11,`$
$`(1\xi )\overline{}u`$ $`=`$ $`1\overline{}(\xi u)\xi \overline{}u=1\xi \overline{}u.`$
For the grades we note that $`(\xi \overline{}f_m^1)e_m=\xi `$ and $`e_m\xi =\xi `$, using (2.5). $``$$``$
###### Proposition 3.3
There is an associator $`\mathrm{\Phi }_{UVW}:(UV)WU(VW)`$ given by
$$\mathrm{\Phi }((\xi \eta )\zeta )=\xi \overline{}\tau (\eta ,\zeta )(\eta \zeta ).$$
Proof First we must check that $`\mathrm{\Phi }`$ preserves the grading. This is just the identity
$$(\xi \eta )\zeta =(\xi \tau (\eta ,\zeta ))(\eta \zeta ).$$
Now we check that the $`G`$ action commutes with the associator. Begin with $`\overline{}u`$
$$\left(\left(\xi \eta \right)\zeta \right)\overline{}u=\left(\xi \overline{}(\eta (\zeta u))\eta \overline{}(\zeta u)\right)\zeta \overline{}u,$$
and apply $`\mathrm{\Phi }`$ to get
$$\xi \overline{}(\eta (\zeta u))\tau (\eta (\zeta u),\zeta u)\left(\eta \overline{}(\zeta u)\zeta \overline{}u\right).$$
Applying $`\mathrm{\Phi }`$ first and then $`\overline{}u`$ we get
$$\left(\xi \overline{}\tau (\eta ,\zeta )\left(\eta \zeta \right)\right)\overline{}u=\xi \overline{}\tau (\eta ,\zeta )(\eta \zeta u)\left(\eta \overline{}(\zeta u)\zeta \overline{}u\right),$$
which is identical to the first expression by the usual identities.
Now we must check that $`\mathrm{\Phi }`$ obeys the pentagon condition, which states that the following two re-bracketings are the same:
$`((VW)Z)U`$ $``$ $`(VW)(ZU)V(W(ZU))`$
$`((VW)Z)U`$ $``$ $`(V(WZ))UV((WZ)U)V(W(ZU))`$
We apply these operations to $`((\xi \eta )\zeta )\upsilon `$, giving
$`((\xi \eta )\zeta )\upsilon `$ $``$ $`(\xi \eta )\overline{}\tau (\zeta ,\upsilon )(\zeta \upsilon )`$
$`=`$ $`(\xi \overline{}(\eta \tau (\zeta ,\upsilon ))\eta \overline{}\tau (\zeta ,\upsilon ))(\zeta \upsilon )`$
$``$ $`\xi \overline{}(\eta \tau (\zeta ,\upsilon ))\tau (\eta \tau (\zeta ,\upsilon ),\zeta \upsilon )`$
$`(\eta \overline{}\tau (\zeta ,\upsilon )(\zeta \upsilon )),`$
$`((\xi \eta )\zeta )\upsilon `$ $``$ $`(\xi \overline{}\tau (\eta ,\zeta )(\eta \zeta ))\upsilon `$
$``$ $`\xi \overline{}\tau (\eta ,\zeta )\tau (\eta \zeta ,\upsilon )((\eta \zeta )\upsilon )`$
$``$ $`\xi \overline{}\tau (\eta ,\zeta )\tau (\eta \zeta ,\upsilon )(\eta \overline{}\tau (\zeta ,\upsilon )(\zeta \upsilon )).`$
These are the same by the usual identities.
We must check the triangle identity, that is the maps $`\mathrm{id}r`$ and $`\mathrm{\Phi }(l\mathrm{id}):VWV(kW)`$ are the same.
$`(\mathrm{id}r)(\xi \eta )`$ $`=`$ $`\xi (1\eta )`$
$`\mathrm{\Phi }(l\mathrm{id})(\xi \eta )`$ $`=`$ $`\mathrm{\Phi }((\xi \overline{}f_m^11)\eta )=\xi \overline{}f_m^1\tau (1,\eta )(1\eta ).`$
These are the same as $`\tau (1,\eta )=\tau (e_m,\eta )=f_m`$, from (2.5).
Finally we check that condition that $`\mathrm{\Phi }`$ is a natural transformation, i.e. that the following diagram commutes,
$$\begin{array}{ccc}(UV)W& \stackrel{\mathrm{\Phi }_{UVW}}{}& U(VW)\\ (\psi \theta )\varphi & & \psi (\theta \varphi )\\ (\stackrel{~}{U}\stackrel{~}{V})\stackrel{~}{W}& \stackrel{\mathrm{\Phi }_{\stackrel{~}{U}\stackrel{~}{V}\stackrel{~}{W}}}{}& \stackrel{~}{U}(\stackrel{~}{V}\stackrel{~}{W})\end{array},$$
where
$`((\psi \theta )\varphi )((\xi \eta )\kappa )`$ $`=`$ $`(\psi (\xi )\theta (\eta ))\varphi (\kappa ),`$
$`(\psi (\theta \varphi ))(\xi (\eta \kappa ))`$ $`=`$ $`\psi (\xi )(\theta (\eta )\varphi (\kappa )).`$
This is simple to check, remembering that the morphisms preserve the grade and action. $``$$``$
## 4 A rigid tensor category
Take a group $`X`$ with subgroup $`G`$, and a set of left coset representatives $`M`$ which contains $`e`$. We suppose that $`(M,)`$ has right inverses, i.e. for every $`sM`$ there is an $`s^RM`$ so that $`ss^R=e`$.
Take a decomposition of an object $`V`$ in $`𝒞`$ according to the grading, i.e. $`V=_{sM}V_s`$, where $`\xi V_s`$ corresponds to $`\xi =s`$. Now take the dual vector space $`V^{}`$, and set
$$V_{s^L}^{}=\{\alpha V^{}:\alpha |_{V_t}=0ts\}.$$
Then $`V^{}=_{sM}V_{s^L}^{}`$, and we define $`\alpha =s^L`$ when $`\alpha V_{s^L}^{}`$. The evaluation map $`\mathrm{ev}:V^{}Vk`$ is defined by $`\mathrm{ev}(\alpha ,\xi )=\alpha (\xi )`$. We have designed the grading on $`V^{}`$ so that this map preserves gradings. Now considering the action $`\overline{}u`$, if we apply evaluation to $`\alpha \overline{}(\xi u)\xi \overline{}u`$ we should get $`\alpha (\xi )\overline{}u=\alpha (\xi )`$. To do this we define $`(\alpha \overline{}(\xi u))(\xi \overline{}u)=\alpha (\xi )`$, or if we put $`\eta =\xi \overline{}u`$,
$$(\alpha \overline{}((\eta u^1)u))(\eta )=\alpha (\eta \overline{}u^1)=(\alpha \overline{}(\eta u^1)^1)(\eta ).$$
If we rearrange this to give $`\alpha v`$ we get the following formula;
$`(\alpha \overline{}v)(\eta )=\alpha (\eta \overline{}\tau (\eta ^L,\eta )^1(\eta ^Lv^1)\tau (\eta ^Lv^1,(\eta ^Lv^1)^R)).`$ (5)
To define the coevaluation map we take a basis $`\{\xi \}`$ of each $`V_s`$, and a corresponding dual basis $`\{\widehat{\xi }\}`$ of each $`V_{s^L}^{}`$, i.e. $`\widehat{\eta }(\xi )=\delta _{\xi ,\eta }`$. Then we put these bases together for all $`sM`$, and define
$$\mathrm{coev}(1)=\underset{\xi \mathrm{basis}}{}\xi \overline{}\tau (\xi ^L,\xi )^1\widehat{\xi }.$$
###### Proposition 4.1
The coevaluation map defined above is a morphism in $`𝒞`$.
Proof First show that each summand in the coevaluation has grade $`e`$. If we put $`s=\xi `$, we have to show that $`(s\tau (s^L,s)^1)s^L=e`$. If we apply $`s`$ to $`(s\tau (s^L,s))s^L`$ we get
$$((s\tau (s^L,s)^1)s^L)s=s(s^Ls)=se=s,$$
so using right division shows that $`(s\tau (s^L,s)^1)s^L=e`$ as required.
It is reasonably easy to see that the map is independent of the choice of basis. If we apply $`\overline{}u`$ to the coevaluation, we get
$$\mathrm{coev}(1)\overline{}u=\underset{\xi }{}\xi \overline{}\tau (\xi ^L,\xi )^1(\xi ^Lu)\widehat{\xi }\overline{}u.$$
Now define a new basis by $`\eta =\xi \overline{}\tau (\xi ^L,\xi )^1(\xi ^Lu)\tau (\xi ^Lu,(\xi ^Lu)^R)`$. We see that $`(\alpha \overline{}u^1)(\xi )=\alpha (\eta )`$, so the dual basis is given by $`\widehat{\xi }=\widehat{\eta }\overline{}u^1`$. Now if we write the coevaluation in terms of the new basis we get
$$\mathrm{coev}(1)=\underset{\eta }{}\eta \overline{}\tau (\eta ^L,\eta )^1\widehat{\eta }.$$
Since $`\widehat{\eta }=\widehat{\xi }u`$, we see that $`\tau (\eta ^L,\eta )=\tau (\xi ^Lu,(\xi ^Lu)^R)`$, so the expressions for $`\mathrm{coev}(1)\overline{}u`$ and $`\mathrm{coev}(1)`$ in the new basis coincide. We conclude that the action is trivial on $`\mathrm{coev}(1)`$ as required. $``$$``$
Now we need to check the consistency of the evaluation, coevaluation and associator. Consider the maps, for a homogenous basis element $`\eta `$:
$`\eta `$ $`\stackrel{\mathrm{coev}I}{}`$ $`{\displaystyle \underset{\xi }{}}(\xi \overline{}\tau (\xi ^L,\xi )^1\widehat{\xi })\eta `$ (6)
$`\stackrel{\mathrm{\Phi }}{}`$ $`{\displaystyle \underset{\xi }{}}\xi \overline{}\tau (\xi ^L,\xi )^1\tau (\xi ^L,\eta )(\widehat{\xi }\eta )`$ (7)
$`\stackrel{I\mathrm{eval}}{}`$ $`{\displaystyle \underset{\xi }{}}\xi \overline{}\tau (\xi ^L,\xi )^1\tau (\xi ^L,\eta )\delta _{\xi ,\eta }=\eta .`$ (8)
$`\widehat{\eta }`$ $`\stackrel{I\mathrm{coev}}{}`$ $`{\displaystyle \underset{\xi }{}}\widehat{\eta }(\xi \overline{}\tau (\xi ^L,\xi )^1\widehat{\xi })`$ (9)
$`\stackrel{\mathrm{\Phi }^1}{}`$ $`{\displaystyle \underset{\xi }{}}(\widehat{\eta }\overline{}\tau (\xi \tau (\xi ^L,\xi )^1,\xi ^L)^1\xi \overline{}\tau (\xi ^L,\xi )^1)\widehat{\xi }.`$ (10)
Now we use the calculation
$`e`$ $`=`$ $`\xi \tau (\xi ^L,\xi )^1\xi ^L=(\xi \tau (\xi ^L,\xi )^1)(\xi \tau (\xi ^L,\xi )^1)\xi ^L`$
$`=`$ $`(\xi \tau (\xi ^L,\xi )^1)\tau (\xi \tau (\xi ^L,\xi )^1,\xi ^L)((\xi \tau (\xi ^L,\xi )^1)\xi ^L),`$
to rewrite the last line of (10) as
$`\stackrel{\mathrm{\Phi }^1}{}`$ $`{\displaystyle \underset{\xi }{}}(\widehat{\eta }\overline{}(\xi \tau (\xi ^L,\xi )^1)\xi \overline{}\tau (\xi ^L,\xi )^1)\widehat{\xi }`$
$`\stackrel{\mathrm{eval}I}{}`$ $`{\displaystyle \underset{\xi }{}}\delta _{\xi ,\eta }\widehat{\xi }=\widehat{\eta }.`$
## 5 An algebra in the tensor category
Take a group $`X`$ with subgroup $`G`$, and a set of left coset representatives $`M`$ which contains $`e`$. We assume that $`(M,)`$ has the left division property, i.e. for all $`t,sM`$ there is a unique solution $`pM`$ to the equation $`sp=t`$.
We can combine the group action and the grading in the definition of $`𝒞`$ by considering a single object $`H`$, a vector space spanned by a basis $`\delta _su`$ for $`sM`$ and $`uG`$. We suppose that $`H`$ is in the category $`𝒞`$, and define a map $`\overline{}:VHV`$ (for $`V`$ any object of $`𝒞`$) by
$$\xi \overline{}(\delta _su)=\delta _{s,\xi }\xi \overline{}u.$$
If this map is to be a morphism in the category we must have $`\xi \delta _su=\xi \overline{}u`$ if $`\xi =s`$, i.e. $`s\delta _su=su`$. This can be solved uniquely for $`\delta _su`$ in $`(M,)`$ by left division. The action of $`vG`$ is given by (using $`a=\delta _su`$)
$`(\delta _su)\overline{}v=\delta _{s(av)}(av)^1uv.`$ (11)
###### Proposition 5.1
The action and grading on $`H`$ are consistent. Further $`\overline{}:VHV`$ is a morphism in $`𝒞`$, for $`V`$ any object of $`𝒞`$.
Proof First we check that $`(\delta _su)\overline{}v=\delta _suv`$. If we set $`b=(\delta _su)\overline{}v`$ then from (11), $`(s(av))b=suv`$, where $`a=\delta _su`$. If we apply $`v`$ to the equation $`s\delta _su=su`$ and use the uniqueness of the result of the left division process, we see $`b=\delta _suv`$.
The grading on $`H`$ was defined so that $`\overline{}`$ preserved the grades, so we only have to check the $`G`$-action. If we set $`a=\delta _su`$ again, then
$$(\xi (\delta _su))\overline{}v=\xi \overline{}(av)(\delta _{s(av)}(av)^1uv),$$
and applying $`\overline{}`$ to this gives
$$\delta _{\xi (av),s(av)}\xi \overline{}uv,$$
which is just $`(\xi \overline{}(\delta _su))\overline{}v`$ as required. $``$$``$
We would now like to give $`H`$ a multiplication so that $`\overline{}`$ becomes an action of the algebra $`H`$. Note that the result is not the usual semi-direct product multiplication.
###### Proposition 5.2
The formula for the product $`\mu `$ for $`H`$ in $`𝒞`$ consistent with action above is
$$(\delta _su)(\delta _tv)=\delta _{t,su}\delta _{s\tau (a,b)}\tau (a,b)^1uv,$$
where $`a=\delta _su`$ and $`b=\delta _tv`$.
Proof We want the following equation to hold, remembering to use $`\mathrm{\Phi }`$ when we change the bracketing:
$`(\xi \overline{}(\delta _su))\overline{}(\delta _tv)=(\xi \overline{}\tau (a,b))\overline{}((\delta _su)(\delta _tv)),`$ (12)
where $`a=\delta _su`$ and $`b=\delta _tv`$. Now
$`(\xi \overline{}(\delta _su))\overline{}(\delta _tv)`$ $`=`$ $`\delta _{s,\xi }(\xi \overline{}u)\overline{}(\delta _tv)`$
$`=`$ $`\delta _{s,\xi }\delta _{t,\xi u}\xi \overline{}uv,`$
and the two sides of (12) agree by definition of the product above. $``$$``$
###### Proposition 5.3
Multiplication $`\mu :HHH`$ is a morphism in $`𝒞`$.
Proof Set $`\eta =\delta _su`$, $`\xi =\delta _tv`$, $`a=\eta `$ and $`b=\xi `$. For the grading, note that by definition $`sa=su`$, $`tb=tv`$ and $`(s\tau (a,b))\eta \xi =suv`$, under the assumption that $`su=t`$. But then $`(sa)b=suv`$, so $`(s\tau (a,b))(ab)=suv`$ and we deduce that $`\eta \xi =ab`$.
Now we check the action:
$$(\eta \xi )\overline{}w=(\delta _{s(a(bw))}(a(bw))^1u(bw))(\delta _{t(bw)}(bw)^1vw),$$
and multiplying these together gives
$$\delta _{t,su}\delta _{s\tau (a,b)(abw)}(abw)^1\tau (a,b)^1uvw,$$
which is the same as $`(\eta \xi )\overline{}w`$. $``$$``$
###### Proposition 5.4
Multiplication $`\mu :HHH`$ is associative in $`𝒞`$. There is an identity $`I`$ for the multiplication and an algebra map $`ϵ:Hk`$ in the category, given by
$$I=\underset{t}{}\delta _te,ϵ(\delta _su)=\delta _{s,e}.$$
In terms of the action of $`H`$ on objects in $`𝒞`$, the identity $`I`$ has the trivial action on all objects. The action of $`hH`$ on the object $`k`$ is just multiplication by $`ϵ(h)`$, and $`ϵ(I)=1`$.
Proof Set $`a=\delta _su`$, $`b=\delta _tv`$ and $`c=\delta _rw`$. Then
$`((\delta _su)(\delta _tv))(\delta _rw)`$ $`=`$ $`\delta _{t,su}(\delta _{s\tau (a,b)}\tau (a,b)^1uv)(\delta _rw)`$
$`=`$ $`\delta _{t,su}\delta _{r,suv}\delta _{s\tau (a,b)\tau (ab,c)}`$
$`\tau (ab,c)^1\tau (a,b)^1uvw,`$
$`((\delta _su)\overline{}\tau (b,c))((\delta _tv)(\delta _rw))`$ $`=`$ $`\delta _{r,tv}(\delta _{s(a\tau (b,c))}(a\tau (b,c))^1u\tau (b,c))`$
$`(\delta _{t\tau (b,c)}\tau (b,c)^1vw)`$
$`=`$ $`\delta _{t,su}\delta _{r,suv}\delta _{s(a\tau (b,c))\tau (a\tau (b,c),bc)}`$
$`\tau (a\tau (b,c),bc)^1(a\tau (b,c))^1uvw,`$
and these are equal by standard identities on $`\tau `$.
For the identity, note that $`I=e`$, which is required as strictly the identity is a morphism $`:kH`$ in the category. The rest is standard. $``$$``$
## 6 A braided tensor category
Take a group $`X`$ with subgroup $`G`$, and a set of left coset representatives $`M`$ which contains $`e`$. We consider a subcategory $`𝒟`$ of $`𝒞`$ with the additional structures of a function $`\overline{}:M\times VV`$ and a $`G`$-grading, written $`|\xi |G`$ for $`\xi `$ in every object $`V`$ in the category $`𝒟`$. We require the following connections between the gradings and actions:
$`|\eta \overline{}u|=(\eta u)^1|\eta |u`$ , $`s\eta =s\overline{}\eta (s|\eta |),`$ (13)
$`\tau (s,\eta )^1(s|\eta |)`$ $`=`$ $`\tau (s\overline{}\eta ,s|\eta |)^1|s\overline{}\eta |.`$ (14)
The operation $`\overline{}`$ is an ‘action’ of $`M`$, which we define to mean that $`t\overline{}:VV`$ is linear for all objects $`V`$ and all $`tM`$, and that
$`p\overline{}(t\overline{}\kappa )`$ $`=`$ $`(p^{}t\overline{}\kappa )\overline{}\tau (p^{}(t|\kappa |),t|\kappa |)^1,`$ (15)
where $`p^{}=p\tau (t\overline{}\kappa ,t|\kappa |)\tau (t,\kappa )^1`$. We also require a cross relation between the two actions,
$`(s\overline{}\eta )\overline{}((s|\eta |)u)`$ $`=`$ $`(s(\eta u))\overline{}(\eta \overline{}u).`$ (16)
The morphisms in the category $`𝒟`$ are linear maps preserving both gradings and both actions.
###### Proposition 6.1
The connections between the gradings and the actions are given by the following factorisations in $`X`$:
$`|s\overline{}\eta |^1s\overline{}\eta `$ $`=`$ $`(s|\eta |)|\eta |^1\eta (s|\eta |)^1,`$
$`|\eta \overline{}u|^1\eta \overline{}u`$ $`=`$ $`u^1|\eta |^1\eta u.`$
Proof Directly from the conditions above. $``$$``$
Now we would like to make $`𝒟`$ into a tensor category. To do this we give the $`G`$-grading and action of $`M`$ on tensor products, and show that the associator is a morphism. We define
$`|\xi \eta |`$ $`=`$ $`\tau (\xi ,\eta )^1|\xi ||\eta |,`$ (17)
$`(s\tau (\eta ,\kappa ))\overline{}(\eta \kappa )`$ $`=`$ $`(s\overline{}\eta )\overline{}\tau (s|\eta |,\kappa )\tau ((s|\eta |)\overline{}\kappa ,s|\eta ||\kappa |)^1(s|\eta |)\overline{}\kappa .`$ (18)
###### Proposition 6.2
The gradings on the tensor product of objects $`VW`$ are given by the following factorisation in $`X`$: $`|\xi \eta |^1\xi \eta =|\eta |^1|\xi |^1\xi \eta .`$
Proof
$$|\xi \eta |^1\xi \eta =|\xi \eta |^1(\xi \eta )=|\xi \eta |^1\tau (\xi ,\eta )^1\xi \eta =|\eta |^1|\xi |^1\xi \eta $$
$``$$``$
###### Proposition 6.3
The gradings on the tensor product are consistent with the actions, as specified in (6.1).
Proof First we check the $`G`$-action. From (6.2), for all $`uG`$,
$`|\xi \overline{}(\eta u)\eta \overline{}u|^1\xi \overline{}(\eta u)\eta \overline{}u`$ $`=`$ $`|\eta \overline{}u|^1|\xi \overline{}(\eta u)|^1\xi \overline{}(\eta u)\eta \overline{}u`$
$`=`$ $`|\eta \overline{}u|^1(\eta u)^1|\xi |^1\xi (\eta u)\eta \overline{}u`$
$`=`$ $`u^1|\eta |^1|\xi |^1\xi \eta u.`$
Now we check the $`M`$ action by considering the grades of $`(s\tau (\eta ,\kappa ))\overline{}(\eta \kappa )`$. We set $`u=\tau (s|\eta |,\kappa )\tau ((s|\eta |)\overline{}\kappa ,s|\eta ||\kappa |)^1`$.
$`|(s\overline{}\eta )\overline{}u(s|\eta |)\overline{}\kappa |^1(s\overline{}\eta )\overline{}u(s|\eta |)\overline{}\kappa `$
$`=|(s|\eta |)\overline{}\kappa |^1|(s\overline{}\eta )\overline{}u|^1(s\overline{}\eta )\overline{}u(s|\eta |)\overline{}\kappa `$
$`=|(s|\eta |)\overline{}\kappa |^1u^1|s\overline{}\eta |^1s\overline{}\eta u(s|\eta |)\overline{}\kappa `$
$`=|(s|\eta |)\overline{}\kappa |^1u^1(s|\eta |)|\eta |^1\eta (s|\eta |)^1u(s|\eta |)\overline{}\kappa .`$
Now use the fact from(14) that $`u|(s|\eta |)\overline{}\kappa |=(s|\eta |)|\kappa |`$ to get
$`|(s\overline{}\eta )\overline{}u(s|\eta |)\overline{}\kappa |^1(s\overline{}\eta )\overline{}u(s|\eta |)\overline{}\kappa `$
$`=((s|\eta |)|\kappa |)^1(s|\eta |)|\eta |^1\eta (s|\eta |)^1((s|\eta |)|\kappa |)|(s|\eta |)\overline{}\kappa |^1(s|\eta |)\overline{}\kappa `$
$`=((s|\eta |)|\kappa |)^1(s|\eta |)|\eta |^1\eta (s|\eta |)^1((s|\eta |)|\kappa |)(s|\eta ||\kappa |)|\kappa |^1\kappa (s|\eta ||\kappa |)^1`$
$`=(s|\eta ||\kappa |)|\kappa |^1|\eta |^1\eta \kappa (s|\eta ||\kappa |)^1`$
$`=(s\tau (\eta ,\kappa )|\eta \kappa |)|\kappa |^1|\eta |^1\eta \kappa (s\tau (\eta ,\kappa )|\eta \kappa |)^1,`$
as required. $``$$``$
###### Proposition 6.4
The function $`\overline{}`$ applied to $`VW`$ satisfies the condition (15) to be an $`M`$-action, and $`\overline{}`$ and $`\overline{}`$ satisfy the cross relation (16) on $`VW`$.
Proof Set $`a=\xi `$ and $`b=\eta `$, and begin with the cross relation, with the following formula derived from the left hand side of (16):
$`((s\tau (a,b))\overline{}(\xi \eta ))\overline{}((s\tau (a,b)|\xi \eta |)u)=((s\overline{}\xi )\overline{}v(s|\xi |)\overline{}\eta )\overline{}((s|\xi ||\eta |)u)`$ (19)
$`=(s\overline{}\xi )\overline{}v((s|\xi |)\overline{}\eta ((s|\xi ||\eta |)u))((s|\xi |)\overline{}\eta )\overline{}((s|\xi ||\eta |)u)`$ (20)
$`=(s\overline{}\xi )\overline{}\tau (s|\xi |,b)((s|\xi |)\overline{}\eta (s|\xi ||\eta |)u)\tau ((s|\xi |)\overline{}\eta ((s|\xi ||\eta |)u),s|\xi ||\eta |u)^1`$ (21)
$`(s|\xi |(bu))\overline{}(\eta \overline{}u)`$ (22)
$`=(s\overline{}\xi )\overline{}\tau (s|\xi |,b)((s|\xi |)bu)\tau ((s|\xi |)\overline{}\eta ((s|\xi ||\eta |)u),s|\xi ||\eta |u)^1`$ (23)
$`(s|\xi |(bu))\overline{}(\eta \overline{}u),`$ (24)
where $`v=\tau (s|\xi |,b)\tau ((s|\xi |)\overline{}\eta ,s|\xi ||\eta |)^1`$. This should be the same as the formula derived from the right hand side of (16):
$`(s\tau (a,b)(abu))\overline{}((\xi \eta )\overline{}u)`$ $`=`$ $`(s\tau (a,b)(abu))\overline{}(\xi \overline{}(bu)\eta \overline{}u)`$ (25)
$`=`$ $`(t\overline{}(\xi \overline{}(bu)))\overline{}w(t|\xi \overline{}(bu)|)\overline{}(\eta \overline{}u),`$ (26)
where we have set
$`t`$ $`=`$ $`s\tau (a,b)(abu)\tau (a(bu),bu)^1=s(a(bu)),`$
$`w`$ $`=`$ $`\tau (t|\xi \overline{}(bu)|,bu)\tau ((t|\xi \overline{}(bu)|)\overline{}(\eta \overline{}u),t|\xi \overline{}(bu)||\eta \overline{}u|)^1.`$
Now we can simplify some pieces of (26):
$`t|\xi \overline{}(bu)|`$ $`=`$ $`t(a(bu))^1|\xi |(bu)=s|\xi |(bu),`$
$`t\overline{}(\xi \overline{}(bu))`$ $`=`$ $`(s(a(bu)))\overline{}(\xi \overline{}(bu))=(s\overline{}\xi )\overline{}((s|\xi |)(bu))`$
$`=`$ $`(s\overline{}\xi )\overline{}\tau (s|\xi |,b)((s|\xi |)bu)\tau (s|\xi |(bu),bu)^1,`$
$`w`$ $`=`$ $`\tau (s|\xi |(bu),bu)\tau ((s|\xi |(bu))\overline{}(\eta \overline{}u),s|\xi |(bu)|\eta \overline{}u|)^1`$
$`=`$ $`\tau (s|\xi |(bu),bu)\tau (((s|\xi |)\overline{}\eta )\overline{}((s|\xi ||\eta |)u),s|\xi ||\eta |u)^1.`$
Substituting these in (26) gives the same result as (24), as required.
Now we check the condition for the $`M`$-action. Begin with
$`(p\tau (s\overline{}\xi v,(s|\xi |)\overline{}\eta ))\overline{}((s\tau (a,b))\overline{}(\xi \eta ))`$ (27)
$`=(p\tau (s\overline{}\xi v,(s|\xi |)\overline{}\eta ))\overline{}((s\overline{}\xi )\overline{}v(s|\xi |)\overline{}\eta )`$ (28)
$`=(p\overline{}((s\overline{}\xi )\overline{}v))\overline{}z(p|(s\overline{}\xi )\overline{}v|)\overline{}((s|\xi |)\overline{}\eta ),`$ (29)
where we have set
$$z=\tau (p|(s\overline{}\xi )\overline{}v|,(s|\xi |)\overline{}\eta )\tau ((p|(s\overline{}\xi )\overline{}v|)\overline{}((s|\xi |)\overline{}\eta ),p|(s\overline{}\xi )\overline{}v||(s|\xi |)\overline{}\eta |)^1.$$
We wish to show that (29) is the same as
$`(p^{}(s\tau (a,b))\overline{}(\xi \eta ))\overline{}\tau (p^{}((s\tau (a,b))|\xi \eta |),s\tau (a,b)|\xi \eta |)^1`$ (30)
$`=(p^{}(s\tau (a,b))\overline{}(\xi \eta ))\overline{}\tau (p^{}(s\tau (a,b))^1(s|\xi ||\eta |),s|\xi ||\eta |)^1,`$ (31)
Where we have set
$`p^{}`$ $`=`$ $`p\tau (s\overline{}\xi v,(s|\xi |)\overline{}\eta )\tau ((s\tau (a,b))\overline{}(\xi \eta ),s\tau (a,b)|\xi \eta |)\tau (s\tau (a,b),ab)^1`$
$`=`$ $`p((s\overline{}\xi v)\tau ((s|\xi |)\overline{}\eta ,s|\xi ||\eta |))`$
$`\tau (s\overline{}\xi v\tau ((s|\xi |)\overline{}\eta ,s|\xi ||\eta |),(s|\xi |)\overline{}\eta s|\xi ||\eta |)\tau (s\tau (a,b),ab)^1`$
$`=`$ $`p((s\overline{}\xi v)v^1\tau (s|\xi |,b))\tau (s\overline{}\xi \tau (s|\xi |,b),(s|\xi |)b)\tau (s\tau (a,b),ab)^1`$
$`=`$ $`p(s\overline{}\xi v)^1(s\overline{}\xi \tau (s|\xi |,b))\tau (s\overline{}\xi \tau (s|\xi |,b),(s|\xi |)b)\tau (s\tau (a,b),ab)^1`$
$`=`$ $`p(s\overline{}\xi v)^1\tau (s\overline{}\xi ,s|\xi |)\tau (s\overline{}\xi (s|\xi |),b)\tau (s\tau (a,b),ab)^1`$
$`=`$ $`p(s\overline{}\xi v)^1\tau (s\overline{}\xi ,s|\xi |)\tau (sa,b)\tau (s\tau (a,b),ab)^1`$
$`=`$ $`p(s\overline{}\xi v)^1\tau (s\overline{}\xi ,s|\xi |)\tau (s,a)^1(s\tau (a,b))`$
If we put $`p^{\prime \prime }=p(s\overline{}\xi v)^1\tau (s\overline{}\xi ,s|\xi |)\tau (s,a)^1`$ then (31) becomes
$`(((p^{\prime \prime }s)\tau (a,b))\overline{}(\xi \eta ))\overline{}\tau (p^{\prime \prime }(s|\xi ||\eta |),s|\xi ||\eta |)^1`$ (32)
$`=((p^{\prime \prime }s\overline{}\xi )\overline{}u^{}((p^{\prime \prime }s)|\xi |)\overline{}\eta )\overline{}\tau (p^{\prime \prime }(s|\xi ||\eta |),s|\xi ||\eta |)^1,`$ (33)
where $`u^{}=\tau ((p^{\prime \prime }s)|\xi |,b)\tau (((p^{\prime \prime }s)|\xi |)\overline{}\eta ,(p^{\prime \prime }s)|\xi ||\eta |)^1`$. We can simplify matters by
$`|(s\overline{}\xi )\overline{}v|`$ $`=`$ $`(s\overline{}\xi v)^1|s\overline{}\xi |v,`$
$`(p|(s\overline{}\xi )\overline{}v|)\overline{}((s|\xi |)\overline{}\eta )`$ $`=`$ $`(q(s|\xi |)\overline{}\eta )\overline{}\tau (q((s|\xi |)|\eta |),s|\xi ||\eta |)^1`$
$`=`$ $`(((p^{\prime \prime }s)|\xi |)\overline{}\eta )\overline{}\tau (q((s|\xi |)|\eta |),s|\xi ||\eta |)^1,`$
$`p\overline{}((s\overline{}\xi )\overline{}v)`$ $`=`$ $`(c\overline{}(s\overline{}\xi ))\overline{}((c|s\overline{}\xi |)v)`$
$`=`$ $`(p^{\prime \prime }s\overline{}\xi )\overline{}\tau (p^{\prime \prime }(s|\xi |),s|\xi |)^1((c|s\overline{}\xi |)v)`$
$`=`$ $`(p^{\prime \prime }s\overline{}\xi )\overline{}\tau (q,s|\xi |)^1(qv),`$
where $`c=p(s\overline{}\xi v)^1`$ and
$`q`$ $`=`$ $`p|(s\overline{}\xi )\overline{}v|\tau ((s|\xi |)\overline{}\eta ,s|\xi ||\eta |)\tau (s|\xi |,b)^1`$
$`=`$ $`p(s\overline{}\xi v)^1|s\overline{}\xi |=p^{\prime \prime }(s|\xi |).`$
Now we can rewrite (29) as
$$((p^{\prime \prime }s)\overline{}\xi )\overline{}\tau (q,s|\xi |)^1(qv)z((p^{\prime \prime }s)|\xi |)\overline{}\eta )\overline{}\tau (q((s|\xi |)|\eta |),s|\xi ||\eta |)^1,$$
so all we have to do now to show that this is equal to (33) is to check that
$`u^{}(((p^{\prime \prime }s)|\xi |)\overline{}\eta \tau (p^{\prime \prime }(s|\xi ||\eta |),s|\xi ||\eta |)^1)`$ $`=`$ $`\tau (q,s|\xi |)^1(qv)z.`$ (34)
To simplify what follows we shall use the substitutions
$`f=((p^{\prime \prime }s)|\xi |)\overline{}\eta `$ , $`n=(s|\xi |)\overline{}\eta `$
$`g=p^{\prime \prime }(s|\xi ||\eta |)`$ , $`h=s|\xi ||\eta |.`$
If we use the result $`p|(s\overline{}\xi )\overline{}v|=p^{\prime \prime }(s|\xi |)v`$, then we can rewrite
$`z`$ $`=`$ $`\tau (p^{\prime \prime }(s|\xi |)v,n)\tau (f\tau (g,h)^1,p^{\prime \prime }(s|\xi |)v|(s|\xi |)\overline{}\eta |)^1`$
$`=`$ $`\tau (qv,n)\tau (f\tau (g,h)^1,g)^1`$
and (34) becomes
$`\tau ((p^{\prime \prime }s)|\xi |,b)\tau (f,gh)^1(f\tau (g,h)^1)`$ $`=`$ $`\tau (q,s|\xi |)^1(qv)\tau (qv,n)\tau (f\tau (g,h)^1,g)^1`$ (35)
$`\tau (q,s|\xi |)\tau (q(s|\xi |),b)`$ $`=`$ $`(qv)\tau (qv,n)\tau ((f\tau (g,h)^1)g,h).`$ (36)
Now we note some equations given by the grades:
$`((f\tau (g,h)^1)g)h`$ $`=`$ $`f(gh)=f((p^{\prime \prime }s)|\xi ||\eta |)=((p^{\prime \prime }s)|\xi |)b=(q(s|\xi |))b,`$
$`((qv)n)h`$ $`=`$ $`(q\tau (s|\xi |,b))(nh)=(q\tau (s|\xi |,b))((s|\xi |)b)=(q(s|\xi |))b.`$
If we multiply (36) on the right by $`(q(s|\xi |))b=((f\tau (g,h)^1)g)h`$ we get
$`\tau (q,s|\xi |)(q(s|\xi |))b`$ $`=`$ $`(qv)\tau (qv,n)((f\tau (g,h)^1)g)h.`$ (37)
But we also have $`(f\tau (g,h)^1)g=(qv)n`$, so
$`q(s|\xi |)b`$ $`=`$ $`(qv)\tau (qv,n)((qv)n)h=(qv)(qv)nh`$ (38)
$`(s|\xi |)b`$ $`=`$ $`vnh=v\tau (n,h)(nh)=\tau (s|\xi |,b)((s|\xi |)b),`$ (39)
which at last verifies (34) and gives the answer! $``$$``$
###### Theorem 6.5
When given the following structures, $`𝒟`$ is a braided tensor category:
The identity object is $`k`$, with trivial gradings and actions.
The associator $`\mathrm{\Phi }`$ and the maps $`l`$ and $`r`$ are defined as for $`𝒞`$.
The braiding $`\mathrm{\Psi }:VWWV`$ is defined by $`\mathrm{\Psi }(\xi \eta )=\xi \overline{}\eta \xi \overline{}|\eta |`$.
Proof The following lemmas. $``$$``$
###### Lemma 6.6
The associator $`\mathrm{\Phi }`$ is a morphism in the category $`𝒟`$.
Proof We begin by checking the $`G`$-grade.
$`|(\xi \eta )\kappa |`$ $`=`$ $`\tau (\xi \eta ,\kappa )^1|\xi \eta ||\kappa |`$
$`=`$ $`\tau (\xi \eta ,\kappa )^1\tau (\xi ,\eta )^1|\xi ||\eta ||\kappa |`$
$`|\xi \overline{}\tau (\eta ,\kappa )(\eta \kappa )|`$ $`=`$ $`\tau (\xi \tau (\eta ,\kappa ),\eta \kappa )^1|\xi \overline{}\tau (\eta ,\kappa )||\eta \kappa |`$
$`=`$ $`\tau (\xi \tau (\eta ,\kappa ),\eta \kappa )^1(\xi \tau (\eta ,\kappa ))^1|\xi ||\eta ||\kappa |.`$
These are equal by the properties of $`\tau `$.
Now we check the $`M`$-action. Set $`\xi =a`$, $`\eta =b`$ and $`\kappa =c`$. To begin,
$`(s\tau (ab,c))\overline{}((\xi \eta )\kappa )`$ $`=`$ $`(s\overline{}(\xi \eta ))\overline{}w(s|\xi \eta |)\overline{}\kappa `$ (40)
$`=`$ $`((t\overline{}\xi )\overline{}u(t|\xi |)\overline{}\eta )\overline{}w(s|\xi \eta |)\overline{}\kappa `$ (41)
$`=`$ $`((t\overline{}\xi )\overline{}u((t|\xi |)\overline{}\eta w)((t|\xi |)\overline{}\eta )\overline{}w)(s|\xi \eta |)\overline{}\kappa ,`$ (42)
where $`s=t\tau (a,b)`$ and
$`w`$ $`=`$ $`\tau (s|\xi \eta |,c)\tau ((s|\xi \eta |)\overline{}\kappa ,s|\xi \eta ||\kappa |)^1,`$
$`u`$ $`=`$ $`\tau (t|\xi |,b)\tau ((t|\xi |)\overline{}\eta ,t|\xi ||\eta |)^1.`$
If we set $`\xi ^{}=\xi \overline{}\tau (b,c)`$ and $`p\tau (a\tau (b,c),bc)=s\tau (ab,c)`$, we would like (42) to equal
$`\mathrm{\Phi }^1((p\tau (a\tau (b,c),bc))\overline{}(\xi ^{}(\eta \kappa )))`$ $`=`$ $`\mathrm{\Phi }^1((p\overline{}\xi ^{})\overline{}v(p|\xi ^{}|)\overline{}(\eta \kappa ))`$ (43)
$`=`$ $`\mathrm{\Phi }^1((p\overline{}\xi ^{})\overline{}v((q\overline{}\eta )\overline{}z(q|\eta |)\overline{}\kappa ))`$ (44)
$`=`$ $`((p\overline{}\xi ^{})\overline{}vx(q\overline{}\eta )\overline{}z)(q|\eta |)\overline{}\kappa ,`$ (45)
where $`q\tau (b,c)=p|\xi ^{}|`$ and
$`v`$ $`=`$ $`\tau (p|\xi ^{}|,bc)\tau ((p|\xi ^{}|)\overline{}(\eta \kappa ),p|\xi ^{}||\eta \kappa |)^1,`$
$`z`$ $`=`$ $`\tau (q|\eta |,c)\tau ((q|\eta |)\overline{}\kappa ,q|\eta ||\kappa |)^1,`$
$`x`$ $`=`$ $`\tau (q\overline{}\eta z,(q|\eta |)\overline{}\kappa )^1.`$
Using the usual identities on $`\tau `$,
$$p|\xi ^{}|=p(a\tau (b,c))^1|\xi |\tau (b,c)=t|\xi |\tau (b,c),$$
or equivalently $`q=t|\xi |`$. Then $`s|\xi \eta |=q|\eta |`$, so the third terms in (42) and (45) are equal. Also we get $`z=w`$, so $`((t|\xi |)\overline{}\eta )\overline{}w=(q\overline{}\eta )\overline{}z`$, and the second terms are equal. Next by using the cross relation on the first term of (45),
$$(p\overline{}\xi ^{})\overline{}vx=(p\overline{}(\xi \overline{}\tau (b,c)))\overline{}vx=(t\overline{}\xi )\overline{}(q\tau (b,c))vx$$
Now we are left with the task of showing that
$`(q\tau (b,c))vx`$ $`=`$ $`u(q\overline{}\eta z).`$ (46)
If we use the formula $`(q\tau (b,c))\overline{}(\eta \kappa )=(q\overline{}\eta z)(q|\eta |)\overline{}\kappa `$ then
$`vx`$ $`=`$ $`\tau (q\tau (b,c),bc)\tau ((q\overline{}\eta z)(q|\eta |)\overline{}\kappa ,q|\eta ||\kappa |)^1\tau (q\overline{}\eta z,(q|\eta |)\overline{}\kappa )^1`$
$`=`$ $`\tau (q\tau (b,c),bc)\tau (q\overline{}\eta z\tau ((q|\eta |)\overline{}\kappa ,q|\eta ||\kappa |),(q|\eta |)c)^1`$
$`((q\overline{}\eta z)\tau ((q|\eta |)\overline{}\kappa ,q|\eta ||\kappa |))^1`$
where we have used $`(q|\eta |)\overline{}\kappa (q|\eta ||\kappa |)=(q|\eta |)c`$. Then
$`(q\tau (b,c))vx`$ $`=`$ $`(q\tau (b,c))\tau (q\tau (b,c),bc)\tau (q\overline{}\eta \tau (q|\eta |,c),(q|\eta |)c)^1`$
$`((q\overline{}\eta z)\tau ((q|\eta |)\overline{}\kappa ,q|\eta ||\kappa |))^1`$
$`=`$ $`\tau (q,b)\tau (qb,c)\tau (q\overline{}\eta \tau (q|\eta |,c),(q|\eta |)c)^1`$
$`((q\overline{}\eta z)\tau ((q|\eta |)\overline{}\kappa ,q|\eta ||\kappa |))^1`$
$`=`$ $`\tau (q,b)\tau (q\overline{}\eta ,q|\eta |)^1(q\overline{}\eta \tau (q|\eta |,c))`$
$`((q\overline{}\eta z)\tau ((q|\eta |)\overline{}\kappa ,q|\eta ||\kappa |))^1`$
$`=`$ $`u(q\overline{}\eta \tau (q|\eta |,c))((q\overline{}\eta \tau (q|\eta |,c))\tau ((q|\eta |)\overline{}\kappa ,q|\eta ||\kappa |)^1)`$
$`=`$ $`u(q\overline{}\eta z),`$
as required, where we have used $`q\overline{}\eta (q|\eta |)=qb`$. $``$$``$
###### Lemma 6.7
The maps $`l_V`$ and $`r_V`$ are morphisms in the category $`𝒟`$.
Proof This is reasonably simple from the definitions, rembering that $`\tau (e,s)=\tau (s,e)=e`$ for all $`sM`$. Only the $`G`$-grade and the $`M`$-action need be checked. $``$$``$
###### Lemma 6.8
The map $`\mathrm{\Psi }:VWWV`$ defined by $`\mathrm{\Psi }(\xi \eta )=\xi \overline{}\eta \xi \overline{}|\eta |`$ is a morphism in the category.
Proof First we check the grades, using (6.1):
$`|\xi \overline{}|\eta ||^1|\xi \overline{}\eta |^1\xi \overline{}\eta \xi \overline{}|\eta |`$ $`=`$ $`|\xi \overline{}|\eta ||^1(\xi |\eta |)|\eta |^1\eta (\xi |\eta |)^1\xi \overline{}|\eta |`$
$`=`$ $`|\xi \overline{}|\eta ||^1\xi \overline{}|\eta ||\eta |^1\eta `$
$`=`$ $`|\eta |^1|\xi |^1\xi |\eta ||\eta |^1\eta =|\eta |^1|\xi |^1\xi \eta .`$
Now we check the $`G`$ action:
$`(\mathrm{\Psi }(\xi \eta ))\overline{}u`$ $`=`$ $`(\xi \overline{}\eta )\overline{}((\xi |\eta |)u)\xi \overline{}|\eta |u,`$
$`\mathrm{\Psi }((\xi \eta )\overline{}u)`$ $`=`$ $`\mathrm{\Psi }(\xi \overline{}(\eta u)\eta \overline{}u)=(\xi (\eta u))\overline{}(\eta \overline{}u)\xi \overline{}(\eta u)|\eta \overline{}u|.`$
The first terms are equal by the cross relation (16), and the second terms are equal by the connections between the grades and the relations (14).
Now we check the $`M`$ action. Set $`a=\xi `$ and $`b=\eta `$.
$`\mathrm{\Psi }((s\tau (a,b))\overline{}(\xi \eta ))`$ $`=`$ $`\mathrm{\Psi }((s\overline{}\xi )\overline{}v(s|\xi |)\overline{}\eta )`$ (47)
$`=`$ $`(s\overline{}\xi v)\overline{}((s|\xi |)\overline{}\eta )(s\overline{}\xi )\overline{}v|(s|\xi |)\overline{}\eta |,`$ (48)
where $`v=\tau (s|\xi |,b)\tau ((s|\xi |)\overline{}\eta ,s|\xi ||\eta |)^1`$. Note that $`v|(s|\xi |)\overline{}\eta |=(s|\xi |)|\eta |`$. Then (48) should be the same as
$`(s\tau (a,b))\overline{}\mathrm{\Psi }(\xi \eta )`$ $`=`$ $`(s\tau (a,b))\overline{}(a\overline{}\eta \xi \overline{}|\eta |)`$ (49)
$`=`$ $`(p\overline{}(a\overline{}\eta ))\overline{}w(p|a\overline{}\eta |)\overline{}(\xi \overline{}|\eta |),`$ (50)
where $`s\tau (a,b)=p\tau (a\overline{}\eta ,a|\eta |)`$ and
$$w=\tau (p|a\overline{}\eta |,a|\eta |)\tau ((p|a\overline{}\eta |)\overline{}(\xi \overline{}|\eta |),p|a\overline{}\eta ||\xi \overline{}|\eta ||)^1.$$
Then $`p|a\overline{}\eta |=s(a|\eta |)`$, so the second term of (48) can be written $`(s\overline{}\xi )\overline{}((s|\xi |)|\eta |)`$, and the second term of (50) as $`(s(a|\eta |))\overline{}(\xi \overline{}|\eta |)`$. These are equal by the cross relation (16). For the first terms,
$`(p\overline{}(a\overline{}\eta ))\overline{}w`$ $`=`$ $`((sa)\overline{}\eta )\overline{}\tau (s(a|\eta |),a|\eta |)^1w`$
$`=`$ $`((sa)\overline{}\eta )\overline{}\tau (s\overline{}\xi ((s|\xi |)|\eta |),s(a|\xi |)|\xi \overline{}|\eta ||)^1,`$
$`(s\overline{}\xi v)\overline{}((s|\xi |)\overline{}\eta )`$ $`=`$ $`((q(s|\xi |))\overline{}\eta )\overline{}\tau (q((s|\xi |)|\eta |),s|\xi ||\eta |)^1,`$
where
$$q=s\overline{}\xi v\tau ((s|\xi |)\overline{}\eta ,s|\xi ||\eta |)\tau (s|\xi |,b)^1=s\overline{}\xi .$$
Now we see that $`q(s|\xi |)=sa`$. $``$$``$
###### Lemma 6.9
The map $`\mathrm{\Psi }`$ satisfies the hexagon identities.
Proof Set $`\xi =a`$, $`\eta =b`$ and $`\kappa =c`$. The following two compositions can be seen to be equal:
$`(\xi \eta )\kappa `$ $`\stackrel{\mathrm{\Phi }}{}`$ $`\xi \overline{}\tau (b,c)(\eta \kappa )`$ (51)
$`\stackrel{\mathrm{\Psi }}{}`$ $`(a\tau (b,c))\overline{}(\eta \kappa )\xi \overline{}\tau (b,c)|\eta \kappa |`$ (52)
$`=`$ $`((a\overline{}\eta )\overline{}u(a|\eta |)\overline{}\kappa )\xi \overline{}|\eta ||\kappa |`$ (53)
$`\stackrel{\mathrm{\Phi }}{}`$ $`(a\overline{}\eta )\overline{}u\tau ((a|\eta |)\overline{}\kappa ,a|\eta ||\kappa )((a|\eta |)\overline{}\kappa \xi \overline{}|\eta ||\kappa |),`$ (54)
$`(\xi \eta )\kappa `$ $`\stackrel{\mathrm{\Psi }I}{}`$ $`(a\overline{}\eta \xi \overline{}|\eta |)\kappa `$ (55)
$`\stackrel{\mathrm{\Phi }}{}`$ $`(a\overline{}\eta )\overline{}\tau (a|\eta |,c)(\xi \overline{}|\eta |\kappa )`$ (56)
$`\stackrel{I\mathrm{\Psi }}{}`$ $`(a\overline{}\eta )\overline{}\tau (a|\eta |,c)((a|\eta |)\overline{}\kappa \xi \overline{}|\eta ||\kappa |),`$ (57)
where $`u=\tau (a|\eta |,c)\tau ((a|\eta |)\overline{}\kappa ,a|\eta ||\kappa |)^1`$.
The hexagon identity for the inverse associator asserts that the following should be equal:
$`\xi (\eta \kappa )`$ $`\stackrel{I\mathrm{\Psi }}{}`$ $`\xi (b\overline{}\kappa \eta \overline{}|\kappa |)`$ (58)
$`\stackrel{\mathrm{\Phi }^1}{}`$ $`(\xi \overline{}\tau (b\overline{}\kappa ,b|\kappa |)^1b\overline{}\kappa )\eta \overline{}|\kappa |`$ (59)
$`\stackrel{\mathrm{\Psi }I}{}`$ $`((a\tau (b\overline{}\kappa ,b|\kappa |)^1)\overline{}(b\overline{}\kappa )\xi \overline{}\tau (b\overline{}\kappa ,b|\kappa |)^1|b\overline{}\kappa |)\eta \overline{}|\kappa |,`$ (60)
$`\xi (\eta \kappa )`$ $`\stackrel{\mathrm{\Phi }^1}{}`$ $`(\xi \overline{}\tau (b,c)^1\eta )\kappa `$ (61)
$`\stackrel{\mathrm{\Psi }}{}`$ $`((a\tau (b,c)^1)b)\overline{}\kappa (\xi \overline{}\tau (b,c)^1(b|\kappa |)\eta \overline{}|\kappa |)`$ (62)
$`\stackrel{\mathrm{\Phi }^1}{}`$ $`((((a\tau (b,c)^1)b)\overline{}\kappa )\overline{}\tau (a\tau (b,c)^1(b|\kappa |),b|\kappa |)^1`$ (64)
$`\xi \overline{}\tau (b,c)^1(b|\kappa |))\eta \overline{}|\kappa |.`$
The third terms in (64) are equal, and the second terms can be seen to be equal by the formula from (14), $`\tau (b,\kappa )^1(b|\kappa |)=\tau (b\overline{}\kappa ,b|\kappa |)^1|b\overline{}\kappa |`$. For the first terms use the condition for an $`M`$ action (15) to get
$$(a\tau (b\overline{}\kappa ,b|\kappa |)^1)\overline{}(b\overline{}\kappa )=(((a\tau (b,c)^1)b)\overline{}\kappa )\overline{}\tau (a\tau (b,c)^1(b|\kappa |),b|\kappa |)^1,$$
as required. $``$$``$
###### Proposition 6.10
The braiding is a natural transformation between the tensor product and its opposite in $`𝒟`$.
Proof The statement just means that the following diagram commutes for all morphisms $`\theta :V\stackrel{~}{V}`$ and $`\varphi :W\stackrel{~}{W}`$:
$$\begin{array}{ccc}VW& \stackrel{\mathrm{\Psi }_{VW}}{}& WV\\ \theta \varphi & & \varphi \theta \\ \stackrel{~}{V}\stackrel{~}{W}& \stackrel{\mathrm{\Psi }_{\stackrel{~}{V}\stackrel{~}{W}}}{}& \stackrel{~}{W}\stackrel{~}{V}\end{array},$$
This is simple to check, remembering that the morphisms preserve the grades and actions. $``$$``$
## 7 A double construction
Take a group $`X`$ with subgroup $`G`$, and a set of left coset representatives $`M`$.
###### Definition 7.1
The set $`Y`$, which is identical to $`X`$, is given a binary operation $``$ defined by
$$(us)(vt)=vust=vu\tau (s,t)(st)u,vG,s,tM.$$
Define the functions $`\stackrel{~}{}:Y\times XY`$ and $`\stackrel{~}{\tau }:Y\times YX`$ by $`y\stackrel{~}{}x=x^1yx`$ and $`\stackrel{~}{\tau }(vt,wp)=\tau (t,p)`$. The function $`\stackrel{~}{}:Y\times XX`$ is defined by
$$vt\stackrel{~}{}wp=v^1wpv^{}=twpt_{}^{}{}_{}{}^{1},\mathrm{where}vt\stackrel{~}{}wp=v^{}t^{}v^{}G,t^{}M.$$
###### Proposition 7.2
The maps $`\stackrel{~}{}`$, $`\stackrel{~}{}`$ and $`\stackrel{~}{\tau }`$ satisfy the six conditions listed in (2.4), with $`(Y,)`$ taking the place of $`(M,)`$, and the group $`X`$ taking the place of $`G`$.
Proof The fourth is immediate. For the sixth, consider
$$(\tau (t,p)^1us\tau (t,p))((vt)(wp))=wvus\tau (t,p)(tp)=wvustp=((us)(vt))(wp).$$
For the second condition we start with
$$v^{}u^{}s^{}t^{}=(u^{}s^{})(v^{}t^{})=((us)(vt))\stackrel{~}{}wp=p^1w^1vustwp,$$
where $`(vt)\stackrel{~}{}wp=p^1w^1vtwp=v^{}t^{}`$. From this we deduce that
$$u^{}s^{}=v_{}^{}{}_{}{}^{1}p^1w^1vustwpt_{}^{}{}_{}{}^{1}=us\stackrel{~}{}v^1wpv^{},$$
as $`v^1wpv^{}=twpt_{}^{}{}_{}{}^{1}`$.
For the third condition, we have $`vt\stackrel{~}{}wpus=v^1wpusv^{\prime \prime }`$, where
$$vt\stackrel{~}{}wpus=v^{\prime \prime }t^{\prime \prime }=v^{}t^{}\stackrel{~}{}usv^{\prime \prime }G,t^{\prime \prime }M.$$
Then $`vt\stackrel{~}{}wpus=v^1wpv^{}v_{}^{}{}_{}{}^{1}usv^{\prime \prime }=(vt\stackrel{~}{}wp)((vt\stackrel{~}{}wp)\stackrel{~}{}us)`$.
For the fifth condition, begin with $`wp\stackrel{~}{}\stackrel{~}{\tau }(us,vt)=\tau (s,t)^1wp\tau (s,t)`$, so
$`\stackrel{~}{\tau }(wp\stackrel{~}{}\stackrel{~}{\tau }(us,vt),usvt)`$ $`=`$ $`\tau (p\tau (s,t),st),`$
$`wp\stackrel{~}{}\tau (s,t)=p\tau (s,t)(p\tau (s,t))^1`$ $`=`$ $`p\tau (s,t),`$
and from these we verify the fifth condition, which is
$$(wp\stackrel{~}{}\stackrel{~}{\tau }(us,vt))\stackrel{~}{\tau }(wp\stackrel{~}{}\stackrel{~}{\tau }(us,vt),usvt)=\stackrel{~}{\tau }(wp,us)\stackrel{~}{\tau }((wp)(us),vt).$$
For the first condition, begin with
$$us\stackrel{~}{}(vt\stackrel{~}{}wp)=us\stackrel{~}{}v^1wpv^{}=u^1v^1wpv^{}u^{},$$
where $`us\stackrel{~}{}v^1wpv^{}=us\stackrel{~}{}(vt\stackrel{~}{}wp)=u^{}s^{}`$. Next
$$(usvt)\stackrel{~}{}wp=vu\tau (s,t)(st)\stackrel{~}{}wp=\tau (s,t)^1u^1v^1wpu^{\prime \prime },$$
where, by the second condition, $`(usvt)\stackrel{~}{}wp=u^{\prime \prime }s^{\prime \prime }=(us\stackrel{~}{}(vt\stackrel{~}{}wp))(vt\stackrel{~}{}wp)=u^{}s^{}v^{}t^{}`$, so we get $`u^{\prime \prime }=v^{}u^{}\tau (s^{},t^{})`$, and hence verify the first condition. $``$$``$
###### Proposition 7.3
The element $`e_y=f_m^1e_m=e`$ is a left identity for $`Y`$ (note it is not in general a right identity), and the operation $`(Y,)`$ has the right division property. The corresponding left inverse is given by the formula $`(vt)^L=v^1t^1`$ (for $`vG`$ and $`tM`$).
Proof To show that $`eY`$ is a left inverse, note that $`eus=ues=us`$ for all $`usY`$ ($`uG`$ and $`sM`$).
To check right division we have to check that there is a unique solution $`wpY`$ ($`wG`$ and $`pM`$) to the equation $`wpus=vt`$. The equation gives $`uw\tau (p,s)(ps)=vt`$, and we can solve the equation $`ps=t`$ to give a unique value of $`p`$. Now $`w=u^1v\tau (p,s)^1`$.
To check the formula for the left identity, $`(v^1t^1)vt=vv^1t^1t=e=e_y`$. $``$$``$
###### Proposition 7.4
If we define $`f_y=e_m=f_mX`$, we see that the conditions in (2.14) are satisfied, using $`f_y`$, $`e_y`$, $`X`$ and $`(Y,)`$ instead of $`f_m`$, $`e_m`$, $`G`$ and $`(M,)`$.
Proof For the first condition, note that $`e_y\stackrel{~}{}x=x^1ex=e_y`$. This implies the second condition, $`e_y\stackrel{~}{}x=e_mxe_m^1`$.
For the fourth condition, $`us\stackrel{~}{}e=us`$, and then $`us\stackrel{~}{}e=ses^1=e`$, the third condition.
For the fifth condition, $`\stackrel{~}{\tau }(e_y,us)=\tau (e_m,s)=f_m=f_y`$.
For the sixth condition, $`us\stackrel{~}{}f_y^1=f_musf_m^1=f_mu(sf_m^1)(sf_m^1)`$, so $`\stackrel{~}{\tau }(us\stackrel{~}{}f_y^1,e_y)=\tau (sf_m^1,e_m)=(sf_m^1)^1`$. Then $`us\stackrel{~}{}f_y^1=sf_m^1(sf_m^1)^1=sf_m^1`$, as required.
For the seventh condition, $`(us\stackrel{~}{}f_y^1)e_m=f_m^1f_musf_m^1e_m=us`$. $``$$``$
Now we return to the case where $`eM`$ for simplicity. We introduce a $`Y`$ valued grading on the objects of $`𝒟`$ by $`\xi =|\xi |^1\xi `$. From our previous calculations we know that $`\eta \overline{}u=\eta \stackrel{~}{}u`$, $`s\overline{}\eta =\eta \stackrel{~}{}(s|\eta |)^1`$ and $`\xi \eta =\xi \eta `$.
###### Proposition 7.5
The map $`\widehat{}:V\times XV`$ defined by $`\xi \widehat{}us=(\xi \overline{}u)\widehat{}s`$ ($`uG`$ and $`sM`$), where
$$\xi \widehat{}s=((s^L|\xi |^1)\overline{}\xi )\overline{}\tau (s^L,s),$$
is a right action of the group $`X`$ on $`V`$, where $`V`$ is any object in $`𝒟`$. Further $`\xi \widehat{}us=\xi \stackrel{~}{}us`$.
Proof First consider the grading;
$$\xi \widehat{}s=(s^L|\xi |^1)\overline{}\xi \stackrel{~}{}\tau (s^L,s)=\xi \stackrel{~}{}s^{L1}\tau (s^L,s)=\xi \stackrel{~}{}s,$$
since $`s^Ls=\tau (s^L,s)`$.
Now we wish to show that
$$(\xi \widehat{}us)\widehat{}vt=\xi \widehat{}usvt=\xi \widehat{}u(sv)\tau (sv,t)((sv)t).$$
It is sufficient to prove this with $`u=e`$, so we need to show
$`((\xi \widehat{}s)\overline{}v)\widehat{}t=(\xi \overline{}(sv)\tau (sv,t))\widehat{}((sv)t).`$ (65)
By using the cross relation (16) we see that, for $`wG`$ and $`pM`$,
$`(\eta \overline{}w)\widehat{}p=((p^Lw^1|\eta |^1)\overline{}\eta )\overline{}(p^Lw^1)^1\tau (p^L,p).`$ (66)
Using (66) we calculate
$`((\xi \widehat{}s)\overline{}v)\widehat{}t`$ $`=`$ $`(((s^L|\xi |^1)\overline{}\xi )\overline{}\tau (s^L,s)v)\widehat{}t`$
$`=`$ $`((t^Lv^1\tau (s^L,s)^1z^1)\overline{}((s^L|\xi |^1)\overline{}\xi ))\overline{}(t^Lv^1\tau (s^L,s)^1)^1\tau (t^L,t),`$
where $`z=|(s^L|\xi |^1)\overline{}\xi |`$. Now, from (15),
$$(t^Lv^1\tau (s^L,s)^1z^1)\overline{}((s^L|\xi |^1)\overline{}\xi )=((p^{}(s^L|\xi |^1))\overline{}\xi )\overline{}\tau (p^{}((s^L|\xi |^1)|\xi |),s^L)^1,$$
where
$`p^{}`$ $`=`$ $`t^Lv^1\tau (s^L,s)^1z^1\tau ((s^L|\xi |^1)\overline{}\xi ,s^L)\tau (s^L|\xi |^1,\xi )^1`$
$`=`$ $`t^Lv^1\tau (s^L,s)^1((s^L|\xi |^1)|\xi |)^1.`$
From this we can calculate
$$(p^{}(s^L|\xi |^1))|\xi |=(t^Lv^1\tau (s^L,s)^1)s^L,$$
so
$`((\xi \widehat{}s)\overline{}v)\widehat{}t`$ $`=`$ $`((((t^Lv^1\tau (s^L,s)^1)s^L)|\xi |^1)\overline{}\xi )`$
$`\overline{}\tau (t^Lv^1\tau (s^L,s)^1,s^L)^1(t^Lv^1\tau (s^L,s)^1)^1\tau (t^L,t),`$
and set this equal to $`((a|\xi |^1)\overline{}\xi )\overline{}y`$. Now consider the right hand side of (65), using (66):
$`(\xi \overline{}(sv)\tau (sv,t))\widehat{}((sv)t)`$ $`=`$ $`((((sv)t)^L\tau (sv,t)^1(sv)^1|\xi |^1)\overline{}\xi )`$
$`\overline{}(((sv)t)^L\tau (sv,t)^1(sv)^1)^1\tau (((sv)t)^L,(sv)t),`$
which we set equal to $`((b|\xi |^1)\overline{}\xi )\overline{}x`$. It is our job to show that $`a=b`$ and $`x=y`$. We use the result, derived from the equations $`(cd)^L\tau (c,d)^1cd=e`$ and $`c^1=\tau (c^L,c)^1c^L`$,
$$(cd)^L\tau (c,d)^1=(d^L\tau (c^L,c)^1)c^L$$
to show that
$$b=((t^L\tau ((sv)^L,sv)^1)(sv)^L)(sv)^1.$$
Now $`((sv)^L(sv))v^1=e`$, so $`((sv)^L(sv)^1)s=e`$, i.e. $`(sv)^L(sv)^1=s^L`$. On the other hand,
$`(sv)^L(sv)^1`$ $`=`$ $`(sv)^L((sv)v^1)=\tau ((sv)^L,sv)v^1\tau ((sv)^L((sv)v^1),s)^1`$
$`=`$ $`\tau ((sv)^L,sv)v^1\tau (s^L,s)^1,`$
and we deduce that $`a=b`$. Next consider
$`b^1x`$ $`=`$ $`(sv)\tau (sv,t)((sv)t)^{L1}\tau (((sv)t)^L,(sv)t)`$
$`=`$ $`(sv)\tau (sv,t)((sv)t)=(sv)(sv)t=svt,`$
$`a^1y`$ $`=`$ $`s^{L1}(t^Lv^1\tau (s^L,s)^1)^1(t^Lv^1\tau (s^L,s)^1)^1\tau (t^L,t)`$
$`=`$ $`s^{L1}\tau (s^L,s)vt^{L1}\tau (t^L,t)=svt,`$
so we deduce that $`x=y`$$``$$``$
###### Proposition 7.6
The $`X`$-action on tensor products in $`𝒟`$ is given by
$$(\xi \eta )\widehat{}x=(\xi \widehat{}(\eta \widehat{}x))\eta \widehat{}x.$$
Proof We begin with
$`(\xi \eta )\widehat{}u`$ $`=`$ $`(\xi \eta )\overline{}u=(\xi \overline{}(\eta u)\eta \overline{}u,`$
so we consider $`\eta \stackrel{~}{}u=\eta ut^1`$ where $`\eta \stackrel{~}{}u=vt`$, i.e. $`t=\eta u`$. Then we see that $`\eta \stackrel{~}{}u=\eta u`$ as required. Continue with
$`(\xi \eta )\widehat{}s\tau (s^L,s)^1`$ $`=`$ $`(\xi \eta )\widehat{}s^{L1}=((s^L|\xi \eta |^1)\overline{}(\xi \eta )`$
$`=`$ $`((s^L|\eta |^1|\xi |^1)\overline{}\xi )\overline{}\tau (s^L|\eta |^1,\eta )\tau ((s^L|\eta |^1)\overline{}\eta ,s^L)^1`$
$`(s^L|\eta |^1)\overline{}\eta .`$
Now we can write
$$(s^L|\eta |^1|\xi |^1)\overline{}\xi =\xi \widehat{}p\tau (p^L,p)^1=\xi \widehat{}p^{L1},$$
where $`p^L=s^L|\eta |^1`$, and rewrite
$$(\xi \eta )\widehat{}s\tau (s^L,s)^1=\xi \widehat{}p^{L1}\tau (p^L,\eta )\tau (p^L\overline{}\eta ,s^L)^1(s^L|\eta |^1)\overline{}\eta .$$
We know that
$$|p^L\overline{}\eta |^1p^L\overline{}\eta =(p^L|\eta |)|\eta |^1\eta (p^L|\eta |)^1=s^L|\eta |^1\eta s^{L1}=|\eta |^1\eta \widehat{}s^{L1},$$
so
$`|\eta |^1\eta \widehat{}s^{L1}`$ $`=`$ $`\eta s^{L1}p^L\overline{}\eta ^1`$
$`=`$ $`\eta (p^L\overline{}\eta s^L)^1\tau (p^L\overline{}\eta ,s^L)^1`$
$`=`$ $`\eta (p^L\overline{}\eta (p^L|\eta |))^1\tau (p^L\overline{}\eta ,s^L)^1`$
$`=`$ $`\eta (p^L\eta )^1\tau (p^L\overline{}\eta ,s^L)^1`$
$`=`$ $`\eta \eta ^1p^{L1}\tau (p^L,\eta )\tau (p^L\overline{}\eta ,s^L)^1,`$
as required $``$$``$
###### Proposition 7.7
The braiding $`\mathrm{\Psi }`$ is given in terms of the $`X`$-action by
$$\mathrm{\Psi }(\xi \eta )=\eta \widehat{}(\xi |\eta |)^1\xi \widehat{}|\eta |,\mathrm{\Psi }^1(\xi ^{}\eta ^{})=\eta ^{}\widehat{}|\xi ^{}\widehat{}\eta ^{}|^1\xi ^{}\widehat{}\eta ^{}.$$
Proof By definition of $`\widehat{}s`$,
$$\eta \widehat{}s\tau (s^L,s)^1=\eta \widehat{}s^{L1}=(s^L|\eta |^1)\overline{}\eta ,$$
so we deduce that $`t\overline{}\eta =\eta \widehat{}(t|\eta |)^1`$. Now
$$\mathrm{\Psi }(\xi \eta )=\xi \overline{}\eta \xi \overline{}|\eta |=\eta \widehat{}(\xi |\eta |)^1\xi \widehat{}|\eta |.\mathrm{}$$
## 8 A bialgebra in the braided category
Take a group $`X`$ with subgroup $`G`$, and a set of left coset representatives $`M`$ which contains $`e`$. We assume that $`(M,)`$ has the left division property, i.e. for all $`t,sM`$ there is a unique solution $`pM`$ to the equation $`sp=t`$.
Introduce a vector space $`D`$ with basis $`\delta _yx`$ for $`yY`$ and $`xX`$. Then we define
$$\xi \widehat{}(\delta _yx)=\delta _{y,\xi }\xi \widehat{}x.$$
We see that $`D`$ is an object of $`𝒟`$ with grade $`y\delta _yx=y\stackrel{~}{}x`$ and action
$$(\delta _yx)\widehat{}z=\delta _{y\stackrel{~}{}(a\stackrel{~}{}z)}(a\stackrel{~}{}z)^1xz,$$
where $`a=\delta _yx`$. Then the associative multiplication $`\mu `$ on $`D`$ consistent with the action is
$$(\delta _yx)(\delta _wz)=\delta _{w,y\stackrel{~}{}x}\delta _{y\stackrel{~}{}\stackrel{~}{\tau }(a,b)}\stackrel{~}{\tau }(a,b)^1xz,$$
where $`b=\delta _wz`$, and is a morphism in $`𝒟`$.
This much we have done before in $`𝒞`$. The additional ingredient we have in $`𝒟`$ is the braiding. We can use the braiding to define a coproduct for $`D`$ which in turn gives the tensor product structure in $`𝒟`$.
###### Proposition 8.1
The coproduct in $`D`$ consistent with the tensor product structure in $`𝒟`$ is
$`\mathrm{\Delta }(\delta _yx)`$ $`=`$ $`{\displaystyle \underset{z,wY:zw=y}{}}\delta _{w\stackrel{~}{}|h||h_{(2)}|^1|h_{(1)}|^1}|h_{(1)}||h_{(2)}||h|^1xh_{(2)}^1`$ (68)
$`\delta _{z\stackrel{~}{}|h||h_{(2)}|^1}|h_{(2)}||h|^1x,`$
where $`h=\delta _yx`$, $`h_{(2)}=|h|^1z^1|h|x^1zx`$, and $`h_{(1)}=|h_{(2)}|hh_{(2)}^1`$.
Proof Begin with
$`(\xi \eta )\widehat{}(\delta _yx)=\delta _{y,\xi \eta }(\xi \eta )\widehat{}x=\delta _{y,\xi \eta }\xi \widehat{}(\eta \stackrel{~}{}x)\eta \widehat{}x.`$ (69)
For $`\mathrm{\Delta }h=h_{(1)}h_{(2)}`$ this should be the same as (using $`h=h_{(1)}h_{(2)}`$)
$`(\xi \eta )(h_{(1)}h_{(2)})`$ $`\stackrel{\mathrm{\Phi }}{}`$ $`\xi \widehat{}\stackrel{~}{\tau }(\eta ,h)(\eta (h_{(1)}h_{(2)}))`$ (70)
$`\stackrel{I\mathrm{\Phi }^1}{}`$ $`\xi \widehat{}\stackrel{~}{\tau }(\eta ,h)((\eta \widehat{}\stackrel{~}{\tau }(h_{(1)},h_{(2)})^1h_{(1)})h_{(2)})`$ (71)
$`\stackrel{I(\mathrm{\Psi }I)}{}`$ $`\xi \widehat{}\stackrel{~}{\tau }(\eta ,h)((h_{(1)}^{}\eta ^{})h_{(2)})`$ (72)
$`\stackrel{I\mathrm{\Phi }}{}`$ $`\xi \widehat{}\stackrel{~}{\tau }(\eta ,h)(h_{(1)}^{}\widehat{}\stackrel{~}{\tau }(\eta ^{},h_{(2)})(\eta ^{}h_{(2)}))`$ (73)
$`\stackrel{\mathrm{\Phi }^1}{}`$ $`(\xi \widehat{}\stackrel{~}{\tau }(\eta ,h)nh_{(1)}^{}\widehat{}\stackrel{~}{\tau }(\eta ^{},h_{(2)}))(\eta ^{}h_{(2)})`$ (74)
$`\stackrel{\widehat{}\widehat{}}{}`$ $`(\xi \widehat{}\stackrel{~}{\tau }(\eta ,h)n)\widehat{}(h_{(1)}^{}\widehat{}\stackrel{~}{\tau }(\eta ^{},h_{(2)}))\eta ^{}\widehat{}h_{(2)},`$ (75)
where
$`n`$ $`=`$ $`\stackrel{~}{\tau }(h_{(1)}^{}\widehat{}\stackrel{~}{\tau }(\eta ^{},h_{(2)}),\eta ^{}h_{(2)})^1`$
$`=`$ $`\stackrel{~}{\tau }(h_{(1)}^{}\stackrel{~}{}\stackrel{~}{\tau }(\eta ^{},h_{(2)}),\eta ^{}h_{(2)})^1,`$
$`h_{(1)}^{}\eta ^{}`$ $`=`$ $`\mathrm{\Psi }(\eta \widehat{}\stackrel{~}{\tau }(h_{(1)},h_{(2)})^1h_{(1)})`$
$`=`$ $`h_{(1)}\widehat{}(\eta \widehat{}\stackrel{~}{\tau }(h_{(1)},h_{(2)})^1|h_{(1)}|)^1\eta \widehat{}\stackrel{~}{\tau }(h_{(1)},h_{(2)})^1|h_{(1)}|.`$
As $`\stackrel{~}{\tau }`$ takes values in $`G`$, we see $`h_{(1)}^{}=h_{(1)}\widehat{}\eta ^{}^1`$. If we set
$`h_{(1)}\widehat{}\eta ^{}^1\stackrel{~}{\tau }(\eta ^{},h_{(2)})`$ $`=`$ $`\delta _{\xi \stackrel{~}{}\stackrel{~}{\tau }(\eta ,h)n}n^1\stackrel{~}{\tau }(\eta ,h)^1(\eta \stackrel{~}{}x)=h^{},`$ (76)
$`h_{(2)}`$ $`=`$ $`\delta _{\eta \stackrel{~}{}\stackrel{~}{\tau }(h_{(1)},h_{(2)})^1|h_{(1)}|}|h_{(1)}|^1\stackrel{~}{\tau }(h_{(1)},h_{(2)})x,`$ (77)
then (69) and (75) agree if $`y=\xi \eta `$. Now define $`vt`$ ($`vG`$ and $`tM`$) by the factorisation $`\eta \stackrel{~}{}x=\eta ^{}h_{(2)}=vt`$, and from this $`t=\eta ^{}h_{(2)}`$ and $`\eta \stackrel{~}{}x=\eta xt^1`$. Also we have $`\eta ^{}^1\tau (\eta ^{},h_{(2)})=h_{(2)}(\eta ^{}h_{(2)})^1=h_{(2)}t^1`$ and $`n=\tau (h^{},t)^1`$. Now $`h^{}\stackrel{~}{}t=h^{}tp^1`$ where $`h^{}\stackrel{~}{}t=h^{}\widehat{}t=up`$ ($`uG`$ and $`pM`$). Then
$`\stackrel{~}{\tau }(\eta ,h)n(h^{}\stackrel{~}{}t)`$ $`=`$ $`\tau (\eta ,h)nh^{}tp^1`$
$`=`$ $`\tau (\eta ,h)(h^{}t)p^1`$
$`=`$ $`\tau (\eta ,h)(\eta h)p^1=\eta hp^1,`$
(using $`h^{}t=\eta h`$ from (75)), so from (77),
$`h_{(1)}\widehat{}h_{(2)}`$ $`=`$ $`h^{}\widehat{}t=\delta _{\xi \stackrel{~}{}\stackrel{~}{\tau }(\eta ,h)n(h^{}\stackrel{~}{}t)}(h^{}\stackrel{~}{}t)^1n^1\stackrel{~}{\tau }(\eta ,h)^1(\eta \stackrel{~}{}x)t`$ (78)
$`=`$ $`\delta _{\xi \stackrel{~}{}\eta hp^1}ph^1x.`$ (79)
Now we calculate, using $`h_{(1)}\stackrel{~}{}h_{(2)}=h^{}\stackrel{~}{}t=up`$,
$$h_{(1)}\widehat{}h_{(2)}\stackrel{~}{}h_{(2)}^1=(h_{(1)}\stackrel{~}{}h_{(2)})^1=(h_{(1)}h_{(2)}p^1)^1,$$
so if we apply $`\widehat{}h_{(2)}^1`$ to (79) we get
$`h_{(1)}`$ $`=`$ $`\delta _{\xi \stackrel{~}{}\eta h(h_{(1)}h_{(2)})^1}h_{(1)}h_{(2)}h^1xh_{(2)}^1`$ (80)
$`=`$ $`\delta _{\xi \stackrel{~}{}\eta \tau (h_{(1)},h_{(2)})^1}\tau (h_{(1)},h_{(2)})xh_{(2)}^1.`$ (81)
From $`h=h_{(1)}h_{(2)}`$ we see that $`\tau (h_{(1)},h_{(2)})^1|h_{(1)}|=|h||h_{(2)}|^1`$, so we can rewrite (77) to give
$`h_{(2)}`$ $`=`$ $`\delta _{\eta \stackrel{~}{}|h||h_{(2)}|^1}|h_{(2)}||h|^1x.`$ (82)
Now we use the definition of $`h_{(2)}`$ on this formula to get
$$(\eta \stackrel{~}{}|h||h_{(2)}|^1)h_{(2)}=|h|^1\eta |h||h_{(2)}|^1h_{(2)}=\eta \stackrel{~}{}x,$$
which we rearrange as $`h_{(2)}=|h|^1\eta ^1|h|x^1\eta x`$. If we set $`\eta =z`$ and $`\xi =z^{}`$, we have the constraint $`\xi \eta =y=|\eta |^1z^{}\eta =z(z^{}\stackrel{~}{}\eta )`$, and if we set $`w=z^{}\stackrel{~}{}\eta `$ then
$`\mathrm{\Delta }(\delta _yx)`$ $`=`$ $`{\displaystyle \underset{z,wY:zw=y}{}}\delta _{w\stackrel{~}{}\tau (h_{(1)},h_{(2)})^1}\tau (h_{(1)},h_{(2)})xh_{(2)}^1`$ (84)
$`\delta _{z\stackrel{~}{}|h||h_{(2)}|^1}|h_{(2)}||h|^1x,`$
where $`h_{(2)}=|h|^1z^1|h|x^1zx`$. We find $`h_{(1)}`$ by solving the equation $`h_{(1)}h_{(2)}=h`$ to get $`h_{(1)}=|h_{(2)}|hh_{(2)}^1`$. Finally we can substitute $`\tau (h_{(1)},h_{(2)})=|h_{(1)}||h_{(2)}||h|^1`$.
$``$$``$
###### Proposition 8.2
The map $`ϵ:Dk`$ defined by $`ϵ(\delta _yx)=\delta _{y,e}`$ is a counit for the coproduct, and $`\mathrm{\Delta }(I)=II`$.
Proof Let $`h=\delta _yx`$. For $`(ϵI)\mathrm{\Delta }`$ begin with
$`(ϵI)\mathrm{\Delta }(\delta _yx)`$ $`=`$ $`{\displaystyle \underset{z,wY:zw=y}{}}\delta _{w,e}\delta _{z\stackrel{~}{}|h||h_{(2)}|^1}|h_{(2)}||h|^1x`$ (85)
$`=`$ $`\delta _{y\stackrel{~}{}|h||h_{(2)}|^1}|h_{(2)}||h|^1x,`$ (86)
where $`h_{(2)}=|h|^1y^1|h|x^1yx`$. But by definition $`yh=|h|^1yh=x^1yx`$, so we deduce that $`h_{(2)}=h`$ and $`h_{(1)}=e`$.
Now for $`(Iϵ)\mathrm{\Delta }`$ begin with
$`(Iϵ)\mathrm{\Delta }(\delta _yx)={\displaystyle \underset{z,wY:zw=y}{}}\delta _{z,e}\delta _{w\stackrel{~}{}|h||h_{(2)}|^1|h_{(1)}|^1}|h_{(1)}||h_{(2)}||h|^1xh_{(2)}^1,`$ (87)
where $`h_{(2)}=|h|^1z^1|h|x^1zx=e`$ and $`h_{(1)}=|h_{(2)}|hh_{(2)}^1=h`$.
The proof that $`\mathrm{\Delta }(I)=II`$ is easy once you notice that for every $`h=\delta _yx`$ term in $`I`$, we have $`x=h=e`$. $``$$``$
As the formula for the coproduct is not very nice, we shall use standard diagramatic arguments to show that $`(D,\mu ,\mathrm{\Delta })`$ is a bialgebra . The pentagon identity means that we do not have to keep track of every re-bracketing done in the course of following a diagram. We just assume that there is a fixed bracketing at the beginning and at the end, and apply the associator as required in between. In Fig 1 we give (in order) the symbols we shall use for the braiding $`\mathrm{\Psi }`$, the action $`\widehat{}:VDV`$, the counit, unit, product and coproduct:
The similarity between the symbols for the action and the counit is not coincidental. The counit is the action of $`D`$ on $`k`$, which is traditionally represented by an invisible line. In Fig 2 we give the definition of the product $`\mu `$ and coproduct $`\mathrm{\Delta }`$ on $`D`$.
Now the proof that $`\mathrm{\Delta }`$ is multiplicative is given as: (Fig 3)
To show that $`\mathrm{\Delta }`$ is coassociative we must first show $`(\mathrm{\Phi }((\xi \eta )\kappa ))\widehat{}h=\mathrm{\Phi }(((\xi \eta )\kappa )\widehat{}h)`$, which is easy enough to check from the definitions. This then means that the following two ways of splitting up the calculation of the action on a triple tensor product are the same: (Fig 4)
## 9 A rigid braided tensor category
We assume the same conditions on $`(M,)`$ as the last section. Note that $`(Y)`$ then has right inverses. The definitions of dual, and the corresponding evaluation and coevaluation maps, considered previously for $`𝒞`$, can also be used in $`𝒟`$. Fig 5(a) and 5(b) show the diagrams we shall use for evaluation and coevaluation. Recall that the morphisms in $`𝒟`$ are required to preserve the actions and gradings. This means that if $`T:VW`$ is a morphism, then we have the picture in Fig 5(c).
We would like to show that $`D`$ is a braided Hopf algebra in the category $`𝒟`$, and in Fig 5(d) we give the definition of antipode $`S:DD`$. Note that this is not the same picture as that in . This is because we are using right actions instead of left actions. We cannot simply reflect the picture in either, as the evaulation and coevaluation morphisms have a definite handedness. In the next proposition we find what the formula for the antipode actually is, and then we go on to check that $`S`$ satisfies the required condition in the axioms of a braided Hopf algebra.
###### Proposition 9.1
Let $`h=\delta _yxD`$. Then $`S(h)=\delta _{y^1|h|h^1}hx^1|h|`$.
Proof Suppose that $`S(h)=\delta _y^{}x^{}`$. Then for $`\xi V`$ we have $`\xi \widehat{}S(h)=\delta _{y^{},\xi }\xi \widehat{}x^{}`$, which by definition is equal to the composition
$`\xi h`$ $`\stackrel{\mathrm{coeval}_V\mathrm{\Psi }^1}{}`$ $`{\displaystyle (\eta \widehat{}\tau (\eta ^L,\eta )^1\widehat{\eta })(h\widehat{}|\xi \widehat{}h|^1\xi \widehat{}h)}`$ (88)
$`\stackrel{\mathrm{\Phi }}{}`$ $`{\displaystyle \eta \widehat{}\tau (\eta ^L,\eta )^1\tau (\eta ^L,\xi h)(\widehat{\eta }(h\widehat{}|\xi \widehat{}h|^1\xi \widehat{}h))}`$ (89)
$`\stackrel{I\mathrm{\Phi }^1}{}`$ $`{\displaystyle \eta ^{}((\alpha h\widehat{}|\xi \widehat{}h|^1)\xi \widehat{}h)}`$ (90)
$`\stackrel{I(\widehat{}I)}{}`$ $`{\displaystyle \eta ^{}(\alpha \widehat{}(h\widehat{}|\xi \widehat{}h|^1)\xi \widehat{}h)}`$ (91)
$`\stackrel{I\mathrm{eval}}{}`$ $`{\displaystyle \eta ^{}(\alpha \widehat{}(h\widehat{}|\xi \widehat{}h|^1))(\xi \widehat{}h)},`$ (92)
where $`\eta ^{}=\eta \widehat{}\tau (\eta ^L,\eta )^1\tau (\eta ^L,\xi h)`$ and $`\alpha =\widehat{\eta }\widehat{}\tau (h\widehat{}|\xi \widehat{}h|^1,\xi \widehat{}h)^1`$. We set $`u=|\xi \widehat{}h|`$, $`s=hu^1`$ and $`t=\xi \widehat{}h`$. We have
$$h\widehat{}u^1=\delta _{y\stackrel{~}{}(h\stackrel{~}{}u^1)}(h\stackrel{~}{}u^1)^1xu^1.$$
Now $`h\stackrel{~}{}u^1=vs`$, so $`h\stackrel{~}{}u^1=hu^1s^1`$, and
$$(\widehat{\eta }\widehat{}\tau (s,t)^1)\widehat{}(h\widehat{}u^1)=\delta _{\eta ^L\stackrel{~}{}\tau (s,t)^1,y\stackrel{~}{}hu^1s^1}\widehat{\eta }\widehat{}\tau (s,t)^1suh^1xu^1.$$
Now we recall that $`\tau (s,t)^1s=(st)t^1`$, where $`st=\xi h`$ as $`\mathrm{\Psi }^1`$ preserves grades. Then
$$(\widehat{\eta }\widehat{}\tau (s,t)^1)\widehat{}(h\widehat{}u^1)=\delta _{y,\eta ^L\stackrel{~}{}(\xi h)t^1uh^1}\widehat{\eta }\widehat{}(\xi h)t^1uh^1xu^1.$$
Now $`u^1t=\xi \widehat{}h=\xi \stackrel{~}{}h`$, so
$`(\widehat{\eta }\widehat{}\tau (s,t)^1)\widehat{}(h\widehat{}u^1)=\delta _{y,\eta ^{LD}\stackrel{~}{}\tau (\xi ,h)^1|\xi |}\widehat{\eta }\widehat{}\tau (\xi ,h)^1|\xi |xu^1.`$ (93)
Set $`w=\tau (\xi ,h)^1|\xi |xu^1`$ and $`z=\xi \stackrel{~}{}h`$. If $`(\widehat{\eta }\widehat{}w)(\xi \widehat{}h)`$ is not zero, then $`\widehat{\eta }\stackrel{~}{}w=z^{LD}`$. Then $`z^{LD}\stackrel{~}{}w^1=\widehat{\eta }=\eta ^{LD}`$, so $`(z^{LD}\stackrel{~}{}w^1)^{RD}=\eta `$. If (93) is not zero, then
$`y`$ $`=`$ $`z^{LD}\stackrel{~}{}w^1\tau (\xi ,h)^1|\xi |`$ (94)
$`=`$ $`(ut^1)\stackrel{~}{}ux^1=(t^1u)\stackrel{~}{}x^1,`$ (95)
$`y^1`$ $`=`$ $`u^1t\stackrel{~}{}x^1=\xi \stackrel{~}{}hx^1,`$ (96)
so we deduce that $`y^{}=\xi =y^1\stackrel{~}{}xh^1`$. By definition $`yh=|h|^1yh=x^1yx`$, so we get $`y^{}=\xi =hx^1y^1xh^1=y^1|h|h^1`$. Now, from (5),
$`(\widehat{\eta }\widehat{}w)(\xi \widehat{}h)`$ $`=`$ $`\widehat{\eta }(\xi \widehat{}h\stackrel{~}{\tau }(z^{LD},z)^1(z^{LD}\stackrel{~}{}w^1)\stackrel{~}{\tau }(z^{LD}\stackrel{~}{}w^1,(z^{LD}\stackrel{~}{}w^1)^{RD}))`$ (97)
$`=`$ $`\widehat{\eta }(\xi \widehat{}h\tau (t^L,t)^1(z^{LD}\stackrel{~}{}w^1)\tau (\eta ^L,\eta ))`$ (98)
$`=`$ $`\widehat{\eta }(\xi \widehat{}h\tau (t^L,t)^1t^Lw^1\eta ^{L1}\tau (\eta ^L,\eta ))`$ (99)
$`=`$ $`\widehat{\eta }(\xi \widehat{}ht^1w^1\eta )`$ (100)
$`=`$ $`\widehat{\eta }(\xi \widehat{}ht^1ux^1|\xi |^1\tau (\xi ,h)\eta )`$ (101)
$`=`$ $`\widehat{\eta }(\xi \widehat{}hx^1y|\xi |^1\tau (\xi ,h)\eta ).`$ (102)
We choose a basis element for the summation to be $`\eta =\xi \widehat{}hx^1y|\xi |^1\tau (\xi ,h)\eta `$, and then
$`\xi \widehat{}x^{}`$ $`=`$ $`\eta ^{}=\xi \widehat{}hx^1y|\xi |^1\tau (\xi ,h)\eta \tau (\eta ^L,\eta )^1\tau (\eta ^L,\xi h)`$ (103)
$`=`$ $`\xi \widehat{}hx^1y|\xi |^1\tau (\xi ,h)\eta ^{L1}\tau (\eta ^L,\xi h)`$ (104)
$`=`$ $`\xi \widehat{}hx^1y|\xi |^1\tau (\xi ,h)(\xi h)(\eta ^L(\xi h))^1`$ (105)
$`=`$ $`\xi \widehat{}hx^1y|\xi |^1\xi h(((\eta ^L\tau (\xi ,h)^1)\xi )h)^1`$ (106)
$`=`$ $`\xi \widehat{}hx^1|h|(((\eta ^L\tau (\xi ,h)^1)\xi )h)^1`$ (107)
Finally, from the top line of (96), we get $`\eta ^{LD}\stackrel{~}{}\tau (\xi ,h)^1=y\stackrel{~}{}|\xi |^1`$. Then if we set $`c=(\eta ^L\tau (\xi ,h)^1)\xi `$, we see that $`vc=(y\stackrel{~}{}|\xi |^1)y^{}`$ (for some $`vG`$). But
$$(y\stackrel{~}{}|\xi |^1)y^{}=(|\xi |y|\xi |^1)\xi =|\xi |^1|\xi |y|\xi |^1\xi =y\xi =|h|h^1,$$
so $`c=h^L`$. We conclude that $`x^{}=hx^1|h|`$. $``$$``$
###### Proposition 9.2
The morphisms $`\mu (IS)\mathrm{\Delta }:DD`$ and $`\mu (SI)\mathrm{\Delta }:DD`$ are both equal to $`1.ϵ:DD`$.
Proof This part of the definition of a braided Hopf algebra can be checked by diagrams. First for $`\mu (IS)\mathrm{\Delta }:DD`$ we have Fig. 6.
Then for $`\mu (SI)\mathrm{\Delta }:DD`$ we have Fig. 7. |
warning/0002/hep-ph0002026.html | ar5iv | text | # Effects of a general set of interactions on neutrino propagation in matterafootnote aafootnote aBased on the article [1] written in collaboration with Sven Bergmann and Yuval Grossman.
## I Introduction
Neutrino physics currently provides the strongest experimental evidence for physics beyond the Standard Model (SM). The atmospheric neutrino anomaly and the solar neutrino problem are best explained by neutrino oscillations. This require massive neutrinos that mix, and hence physics beyond the SM. When neutrinos propagate in matter, the physics of neutrino oscillations can be very different from the case of vacuum propagation. This is because coherent interactions with the background give to the neutrino an “index of refraction” that depends on the type of background and on the neutrino flavor. For example, in normal matter only electron neutrinos have SM charged current interactions, and thus the effective $`\nu _e`$ mass is enhanced with respect to the other flavors. This allows for the possibility of level crossing between different neutrino eigenstates in matter, and can result in a significant amplification of neutrino oscillations. This is known as the MSW effect. For light sterile neutrinos also neutral current interactions are important. Finally, in a polarized medium the neutrino effective mass depends also on the average polarization of the background and on the angle between the neutrino momentum and the polarization vector.
New physics models that imply massive neutrinos often predict also new neutrino interactions, that can significantly modify the SM picture. For example, non-universal interactions may give rise to matter effects that distinguish between muon and tau neutrinos. Lepton flavor violating interactions can induce an effective mixing in matter, allowing for a resonant conversion even in the absence of vacuum mixing. The two effects combined together can induce neutrino flavor transitions even for massless neutrinos. Most of the discussions of these non standard effects assume new interactions just of vector and axial vector types. However, recently the possible effects of a much more general set of interactions have been analyzed. In this talk we discuss the main results of this investigation.
## II Neutrino propagation in matter with general interactions
Our aim is to study neutrino propagation in matter in the presence of the most general pointlike and Lorentz invariant four-fermion interaction with the background fermions ($`f=e,p,n,\nu `$). We assume an interaction Hamiltonian of the form
$$_{\mathrm{int}}=\frac{G_F}{\sqrt{2}}\underset{a}{}(\overline{\nu }\mathrm{\Gamma }^a\nu )\left[\overline{\psi }_f\mathrm{\Gamma }_a(g_a+g_a^{}\gamma ^5)\psi _f\right]+\mathrm{h}.\mathrm{c}.,$$
(1)
where $`\mathrm{\Gamma }^a=\{I,\gamma ^5,\gamma ^\mu ,\gamma ^\mu \gamma ^5,\sigma ^{\mu \nu }\}`$, $`\sigma ^{\mu \nu }=\frac{i}{2}[\gamma ^\mu ,\gamma ^\nu ]`$ and $`a=\{S,P,V,A,T\}`$. The Fermi constant $`G_F`$ has been factored out so that all the couplings are dimensionless. In general, $`\nu `$ is a vector of the different neutrino types, and $`g_a,g_a^{}`$ are 10 matrices in the space of neutrino flavors that describe the coupling strengths. Note that new interactions can include both flavor diagonal and off-diagonal couplings. To derive the equation of motion for the neutrino propagation in matter we first average the effective interactions over the background fermions. We select only coherent transitions, which leave the many-fermion background system in the same state, since incoherent effects become negligible after averaging. In particular, while we allow for neutrino spin-flips, we require that the background fermions do not change their spin. Accordingly, we introduce matrix elements of the fermion currents between initial and final states with the same quantum numbers
$$_a^ff,𝒑,𝝀|\overline{\psi }_f\mathrm{\Gamma }_a(g_a+g_a^{}\gamma ^5)\psi _f|f,𝒑,𝝀,$$
(2)
where $`𝒑`$ and $`𝝀`$ denote respectively the momentum and polarization vectors of the background fermion $`f`$. The expectation value of $`_a^f`$, averaged over the fermion distribution $`\rho _f(𝒑,𝝀)`$ reads
$$V_a^f=\frac{G_F}{\sqrt{2}}\underset{𝝀}{}\frac{d^3p}{(2\pi )^3}\rho _f(𝒑,𝝀)_a^f.$$
(3)
The effect of the medium on the neutrino propagation in the presence of the general interactions (1) is then described by the interaction Lagrangian
$$_{int}=\underset{a,f}{}(\overline{\nu }\mathrm{\Gamma }^a\nu )V_a^f.$$
(4)
The computation of the various $`_a^f`$ is straightforward. After performing the contractions $`\mathrm{\Gamma }^aV_a^f`$ in (4) we obtain
$`\mathrm{\Sigma }^{SP}`$ $``$ $`\left[V^S+V^P\gamma ^5\right]={\displaystyle \frac{G_F}{\sqrt{2}}}n_f{\displaystyle \frac{m_f}{E_f}}\left(g_S+g_P^{}\gamma ^5\right)`$ (5)
$`\mathrm{\Sigma }^{VA}`$ $``$ $`\gamma ^\mu \left[V_\mu ^V+V_\mu ^A\gamma ^5\right]`$ (6)
$`=`$ $`{\displaystyle \frac{G_F}{\sqrt{2}}}n_f\left[{\displaystyle \frac{p/}{E_f}}\left(g_V+g_A^{}\gamma ^5\right)+m_f{\displaystyle \frac{s/}{E_f}}\left(g_V^{}+g_A\gamma ^5\right)\right]`$ (7)
$`\mathrm{\Sigma }^T`$ $``$ $`\varsigma ^i\left[V_i^B+iV_i^E\gamma ^5\right]={\displaystyle \frac{G_F}{\sqrt{2}}}n_f{\displaystyle \frac{[s/,p/]}{E_f}}\left(g_T^{}+g_T\gamma ^5\right),`$ (8)
where $`\varsigma ^i\text{diag}(\sigma ^i,\sigma ^i)`$. The spin vector $`s`$ satisfies $`s^2=1`$ and $`s_\mu p^\mu =0`$ (the explicit expression can be found in ). We have also introduced
$$n_f=\underset{𝝀}{}\frac{d^3p}{(2\pi )^3}\rho _f(𝒑,𝝀),x=\frac{1}{n_f}\underset{𝝀}{}\frac{d^3p}{(2\pi )^3}\rho _f(𝒑,𝝀)x(𝒑,𝝀)$$
(9)
to denote, respectively, the number density $`n_f`$ of the fermion $`f`$ and the average of some function $`x(𝒑,𝝀)`$ over the fermion distribution. In (8) we have decomposed the tensor term $`V_{\mu \nu }^T`$ in analogy to the electro-magnetic field tensor $`F_{\mu \nu }`$, as $`V_i^B=ϵ_{ijk}V_{jk}^T`$ and $`V_i^E=2V_{0i}^T`$. Note that the second equality in (8) makes apparent that the tensor interaction can contribute only in the presence of a polarized background.
The equation of motion for the neutrino propagation can be deduced from the Lagrangian
$$=_{free}+_{int}=\overline{\nu }(i/m\mathrm{\Sigma })\nu $$
(10)
where the matrix of the potentials $`\mathrm{\Sigma }\mathrm{\Sigma }^{SP}+\mathrm{\Sigma }^{VA}+\mathrm{\Sigma }^T`$ depends on the background density and polarization, and in general will vary along the neutrino propagation path. In the general case both $`\mathrm{\Sigma }`$ and $`m`$ are matrices in the space of neutrino types. In the chiral basis the interaction part in (10) reads
$$_{int}=\overline{\nu }\mathrm{\Sigma }\nu =\left(\begin{array}{c}\nu _L^{}\\ \nu _R^{}\end{array}\right)^T\left(\begin{array}{cc}V_\mu ^{LL}\overline{\sigma }^\mu & V_\mu ^{LR}\sigma ^\mu \\ V_\mu ^{RL}\sigma ^\mu & V_\mu ^{RR}\sigma ^\mu \end{array}\right)\left(\begin{array}{c}\nu _L\\ \nu _R\end{array}\right),$$
(11)
where $`\sigma ^\mu (\overline{\sigma }^\mu )=(\sigma ^0,()\sigma ^i)`$ with $`\sigma ^0=I`$, and
$`V_\mu ^{LL}`$ $``$ $`V_\mu ^VV_\mu ^A,V_0^{RL}V^SV^P,V_i^{RL}V_i^BiV_i^E,`$ (12)
$`V_\mu ^{RR}`$ $``$ $`V_\mu ^V+V_\mu ^A,V_0^{LR}V^S+V^P,V_i^{LR}V_i^B+iV_i^E.`$ (13)
It is apparent that the (axial)vector potentials appearing in the diagonal entries in (11) couple neutrinos of the same chirality, while the (off-diagonal) (pseudo)scalar and tensor potentials couple neutrinos of opposite chirality. The equations of motion for neutrinos and antineutrinos derived from (10) read
$$\gamma _0(k/m\mathrm{\Sigma })u=0,\gamma _0(k/+m+\mathrm{\Sigma })v=0.$$
(14)
Note that the signs of $`m`$ and $`\mathrm{\Sigma }`$ are opposite for the antineutrinos. The dispersion relations for the neutrino propagation are given by the solutions of
$$\mathrm{det}[𝒪]=\mathrm{det}[\gamma _0(k/m\mathrm{\Sigma })]=0.$$
(15)
Let us take the neutrino momentum $`𝒌=k\widehat{𝒛}`$ along the $`z`$-axis. Assuming that $`V^{V,A,T},mE`$ and neglecting terms of $`𝒪(m/E)`$ the relevant terms in (11) are $`V_{0,3}^{LL}`$, $`V_{0,3}^{RR}`$ and the tensor components $`V_{1,2}^{LR}`$, $`V_{1,2}^{RL}`$ transverse with respect to the neutrino propagation direction. Solving the determinant equation (15) yields the neutrino energies
$$E_\pm =k+\frac{m^2}{2k}+\frac{1}{2}\left[V_{03}^{LL}+V_{03}^{RR}\pm \sqrt{\left(V_{03}^{LL}V_{03}^{RR}\right)^2+4V_+^{LR}V_{}^{RL}}\right],$$
(16)
where $`V_{0\pm 3}V_0\pm V_3`$ and $`V_\pm V_1\pm iV_2`$. In (16) the plus (minus) sign refers to neutrinos that are mainly left(right)-handed states. Eliminating the two helicity suppressed states from the equations of motion we obtain a Schrödinger-like equation that governs the neutrino propagation:
$$i\frac{d}{dt}\left(\begin{array}{c}\nu _L\\ \nu _R\end{array}\right)=_\nu \left(\begin{array}{c}\nu _L\\ \nu _R\end{array}\right)\text{with}_\nu =k+\frac{m^2}{2k}+\left(\begin{array}{cc}V_{03}^{LL}& V_+^{LR}\\ V_{}^{RL}& V_{03}^{RR}\end{array}\right).$$
(17)
The two eigenvalues of the effective Hamiltonian $`_\nu `$ are the solutions (16) of the determinant equation (15). The results for the antineutrinos can be obtained from (17) and (16) by changing the sign of the potentials ($`VV`$). In the case of more than one neutrino flavor (16) is a matrix equation in the space of the neutrino types. In the one flavor case, the energy gap between the two states is
$$\mathrm{\Delta }E_\nu =\sqrt{\left(V_{03}^{LL}V_{03}^{RR}\right)^2+4V_+^{LR}V_{}^{RL}}.$$
(18)
In the limit of vanishing tensor interaction $`(V_T=0)`$ $`\nu _L`$ decouples from $`\nu _R`$. Setting also $`V^{RR}=0`$ and $`V^{LL}`$ equal to the SM charged current and neutral current interactions, we recover the SM case, with decoupled non-interacting right-handed states.
## III Implications and Discussion
Setting $`g_V`$=$`g_A^{}`$=$`g_A`$=$`g_V^{}`$=$`\mathrm{\hspace{0.17em}1}`$ in (7) and $`\mathrm{\Sigma }^{SP}=\mathrm{\Sigma }^T=0`$ yields the SM result for the potential felt by an electron neutrino propagating in an electron background. Introducing a unit vector in the direction of the neutrino momentum $`\widehat{𝒌}=𝒌/|𝒌|`$ and using the explicit expression for the spin vector $`s`$ we reproduce the result
$$V_{\nu ,\overline{\nu }}^{SM}=\pm \sqrt{2}G_Fn_e\left[1\frac{(𝒑+m_e𝝀)\widehat{𝒌}𝒑𝝀}{E_e}+\frac{(\widehat{𝒌}𝒑)(𝒑𝝀)}{E_e(m_e+E_e)}\right]$$
(19)
where the plus-sign refers to neutrinos and the minus-sign to antineutrinos.
However, our main result is that in the presence of a neutrino tensor interaction with the background fermions, the neutrino can undergo spin-flip. This effect is similar to the spin-precession induced by a transverse magnetic field $`B_{}`$ that couples to the neutrino magnetic dipole moment $`\mu _\nu `$. In fact, if we substitute in (17) the off-diagonal term $`V_\pm ^{LR}`$ by $`\mu _\nu B_{}`$ we obtain the equation of motion for a neutrino that propagates in a magnetic field. Of course the two scenarios originate from different physics, however, formally they can be treated in the same way.
In the simplest one generation case, a left-handed neutrino produced at $`t=0`$ and propagating for a time $`t`$ in a constant medium will be converted into a right-handed neutrino with a probability $`P_\nu ^{LR}(t)=\mathrm{sin}^22\theta \mathrm{sin}^2\left(\mathrm{\Delta }E_\nu t/2\right)`$ where the effective mixing angle is given by
$$\mathrm{sin}^22\theta =\frac{|2V_+^{LR}|^2}{(\mathrm{\Delta }E_\nu )^2},$$
(20)
and $`\mathrm{\Delta }E_\nu `$ is given in (18). In the case of more than one neutrino flavor, propagation in a medium with changing density can lead to resonance effects in complete analogy to the magnetic field induced resonant spin-flip.
Now, let us discuss shortly the results for different types of background matter. First consider a medium where the average momentum of the background fermions vanishes: $`𝒑=0`$ (e.g. when the momentum distribution is isotropic). The tensor component determining the effective mixing (20) is given by
$$|V_+^{LR}|=\sqrt{2}G_Fn_f\sqrt{|g_T|^2+|g_T^{}|^2}\lambda _{}\left(\mathrm{sin}^2\vartheta +\frac{m_f}{E_f}\mathrm{cos}^2\vartheta \right),$$
(21)
where $`\vartheta `$ denotes the angle between the momentum and the transverse polarization of the background fermion, and $`\lambda _{}=\sqrt{\lambda _1^2+\lambda _2^2}`$. Note that $`|V_+^{LR}|`$ vanishes if the neutrino propagates along the direction of the average background polarization ($`\lambda _{}=0`$). For a non-relativistic background ($`E_fm_fp_i`$) this yields $`|V_+^{LR}|=\sqrt{2}G_Fn_f\sqrt{|g_T|^2+|g_T^{}|^2}\lambda _{}`$ while in the ultra-relativistic limit we find $`|V_+^{LR}|\lambda _{}\mathrm{sin}^2\vartheta `$. Finally, for a degenerate background in the presence of a magnetic field, only the fermions in the lowest Landau level contribute to the polarization, with the spin oriented antiparallel to the momentum. In this case the background is not isotropic, and one obtains
$$|V_+^{LR}|=\sqrt{2}G_Fn_f\sqrt{|g_T|^2+|g_T^{}|^2}\lambda _{}\frac{m_f}{E_f},$$
(22)
which vanishes in the ultra-relativistic limit.
It is interesting to note that tensor interactions could result from neutrino scalar couplings after Fierz rearrangement. Consider the tree level Lagrangian
$$_{\mathrm{tree}}=\lambda _\varphi \varphi (\overline{L_L}e_R)+\lambda _\varphi ^{}\stackrel{~}{\varphi }(\overline{L_L}\nu _R)+\mathrm{h}.\mathrm{c}.,$$
(23)
involving a right-handed neutrino singlet ($`\nu _R`$) and a doublet scalar field $`\varphi `$ with mass $`m_\varphi `$ and couplings $`\lambda _\varphi ,\lambda _\varphi ^{}`$ to the lepton fields. The resulting set of low energy effective interactions contains the following terms:
$$_{\mathrm{int}}^\varphi =\frac{\lambda _\varphi ^{}\lambda _\varphi }{m_\varphi ^2}\left[\frac{1}{2}(\overline{\nu _R}\nu _L)(\overline{e_R}e_L)+\frac{1}{8}(\overline{\nu _R}\sigma _{\mu \nu }\nu _L)(\overline{e_R}\sigma ^{\mu \nu }e_L)\right],$$
(24)
implying $`g_T\lambda _\varphi ^{}\lambda _\varphi /m_\varphi ^2`$. When different scalar fields mix, operators of this kind can be generated also in supersymmetric models without $`R`$-parity.
Let us now address the issue whether the new tensor term could be relevant for real physical systems, like the sun or a galactic supernova. From eqs. (21)–(22) it follows that, with respect to the SM vector potential, the effective tensor potential is suppressed by a factor
$$ϵ\left|\frac{V_+^{LR}}{V_0^{LL}}\right|<\sqrt{|g_T|^2+|g_T^{}|^2}\lambda _{}.$$
(25)
New physics effects can be relevant to neutrino oscillations only if $`g_T,g_T^{}`$ and $`\lambda _{}`$ are large enough to affect sizably the results obtained within the SM. In particular $`ϵ`$ should satisfy the lower limits: $`ϵ_{sun}>10^2`$ and $`ϵ_{SN}>10^4.`$ The excellent agreement between the SM predictions and various experimental results, suggests that $`g_T`$ and $`g_T^{}`$ are small, probably not exceeding the few percent level. However, the tiny values of the average polarization is by far the most important suppression factor. In the solar interior, the magnetic field can be at most of the order of several kG. This results in a very small electron polarization $`\lambda _e10^8`$ and quite likely neutrino propagation in the sun cannot be affected by the new tensor interaction. For a proto-neutron star in the early cooling phase, soon after the supernova explosion, the magnetic field strength can reach very large values. However, the temperature is also large, thus suppressing the induced polarization. For a magnetic field $`B10^{13}`$ G, it was estimated $`\lambda _e10^4`$ and for the nucleon polarization $`\lambda _{p,n}10^5`$. Thus, if $`B<10^{13}`$ G the propagation of supernova neutrinos would not be affected. However, it was pointed out that collision effects could increase the production of right-handed states and thus enhance the effects of the tensor interaction. Also, the value of the proto-neutron star magnetic field is poorly known. It has been proposed that at early times it could be as large as $`10^{16}`$G. This would imply an enhancement of the polarization of about three orders of magnitude, opening the possibility of observing these effects.
Finally, let us note that since the presence of right-handed neutrinos implies in general a non-vanishing magnetic moment, the effect of the tensor interaction will be accompanied by similar effects due to the neutrino magnetic moment coupled to the strong magnetic field. Clearly, in this case both effects have to be taken into account simultaneously. |
warning/0002/astro-ph0002286.html | ar5iv | text | # Black–Hole Transients and the Eddington Limit
## 1 Introduction
It is now well understood that the accretion discs in low–mass X–ray binaries (LMXBs) are strongly irradiated by the central X–rays, and that this has a decisive effect on their thermal stability (van Paradijs, 1996; King, Kolb & Burderi, 1996). Irradiation stabilizes LMXB discs compared with the otherwise similar ones in cataclysmic variables (CVs) by removing their hydrogen ionization zones. In CVs this instability causes dwarf nova outbursts, and in LMXBs it produces transient outbursts rather than persistent accretion. The irradiation effect appears to be weaker if the accretor is a black hole rather than a neutron star, possibly because of the lack of a hard surface (King, Kolb & Szuszkiewicz, 1997). The result is that neutron–star LMXBs with short ($``$ hours) orbital periods tend to be persistent, while similar black–hole binaries are largely transient. Both types of LMXBs must be transient at sufficiently long orbital periods, since a long period implies a large disc, so that a large X–ray luminosity would be needed to keep the disc edge ionized and thus suppress outbursts. We can write this stability requirement as
$$\dot{M}_{\mathrm{crit}}^{\mathrm{irr}}R_\mathrm{d}^2P^{4/3},$$
(1)
where $`\dot{M}_{\mathrm{crit}}^{\mathrm{irr}}`$ is the minimum central accretion rate required to keep the disc stable, $`R_\mathrm{d}`$ is the outer disc radius, and $`P`$ is the orbital period, and we have used Kepler’s law. Thus for large $`P`$, $`\dot{M}_{\mathrm{crit}}^{\mathrm{irr}}`$ will rise above any likely steady accretion rate, making long–period systems transient. This simple prediction (King, Frank, Kolb & Ritter, 1997) seems to be borne out by the available evidence.
## 2 The Critical Accretion Rate
The precise coefficient in (1) depends on uncertainties in the vertical disc structure (see the discussion in Dubus, Lasota, Hameury & Charles, 1999). Here I adopt the form derived by King, Kolb & Szuszkiewicz (1997). They argued that for a steady black–hole accretor, the central irradiating source is likely to be the inner disc rather than a solid spherical surface, as for a steady neutron–star accretor. (Note that during an outburst of a transient black–hole system such a spherical source may be present, as the central accretor may develop a corona.) For a small source at the centre of the disc and lying in its plane, the irradiation temperature $`T_{\mathrm{irr}}(R)`$ is given by
$$T_{\mathrm{irr}}(R)^4=\frac{\eta \dot{M}c^2(1\beta )}{4\pi \sigma R^2}\left(\frac{H}{R}\right)^2\left(\frac{\mathrm{d}\mathrm{ln}H}{\mathrm{d}\mathrm{ln}R}1\right)$$
(2)
(Fukue, 1992). Here $`\eta `$ is the efficiency of rest–mass energy conversion into X–ray heating, $`\beta `$ is the X–ray albedo, and $`H(R)`$ is the local disc scale height. The minimum accretion rate required to keep the disc in the high state is given by setting $`T_{\mathrm{irr}}(R)=T_H`$, where is $`T_H`$ is the hydrogen ionization temperature. Since $`T`$ always decreases with $`R`$, the global minimum value $`\dot{M}_{\mathrm{crit}}^{\mathrm{irr}}`$ is given by conditions at the outer edge $`R_\mathrm{d}`$ of the disc. For the parametrization adopted by King, Kolb & Szuszkiewicz (1997), and $`\eta =0.2`$, this leads to
$$\dot{M}_{\mathrm{crit}}^{\mathrm{irr}}(R)=2.86\times 10^{11}m_1^{5/6}m_2^{1/2}f_{0.7}^2gP_h^{4/3}\mathrm{M}_{}\mathrm{yr}^1,$$
(3)
where $`f_{0.7}`$ is the disc filling fraction $`f`$ (the ratio of $`R_\mathrm{d}`$ to the accretor’s Roche lobe) in units of 0.7; $`m_1,m_2`$ are the accretor and companion star mass in $`\mathrm{M}_{}`$; and
$`g=`$
$`\left({\displaystyle \frac{1\beta }{0.1}}\right)^1\left({\displaystyle \frac{H}{0.2R}}\right)^2\left({\displaystyle \frac{2/7}{\mathrm{d}\mathrm{ln}H/\mathrm{d}\mathrm{ln}R1}}\right)\left({\displaystyle \frac{T_H}{6500\mathrm{K}}}\right)^4`$ . (4)
Equation (3) is the same as eqn (12) of King, Kolb & Szuszkiewicz (1997) apart from the factors $`f_{0.7}^2g`$, there taken as unity. All of the uncertainties over disc thickness, warping, albedo etc are lumped into the quantity $`g`$. With $`gf_{0.7}1`$, equation (3) appears to be largely successful in predicting that systems with reasonably massive ($`57\mathrm{M}_{}`$) black holes and main–sequence companions should be transient. By contrast, neutron star systems with main–sequence companions should be persistent, as the index of the ratio $`H/R`$ in (2) is unity, implying more efficient disc irradiation. (Equation (3) also implies that lower–mass black hole systems might be persistent.) These results suggest that the quantity $`g`$ appearing in (3) cannot be too far from unity.
## 3 The Eddington Limit
Here I concentrate an another aspect of eq. (3) which does not seem to have received much attention. Namely, for large enough $`P`$, $`\dot{M}_{\mathrm{crit}}^{\mathrm{irr}}`$ must exceed the Eddington accretion rate
$$\dot{M}_{\mathrm{Edd}}1\times 10^8m_1\mathrm{M}_{}\mathrm{yr}^1.$$
(5)
The obvious consequence of eqs. (3, 5) is that for sufficiently long orbital periods irradiation will be unable to suppress outbursts, as the required central luminosity exceeds the Eddington limit, and the system presumably cannot be both super–Eddington and persistent. Note that this conclusion holds whatever the actual value of the mass transfer rate in the particular binary happens to be. Thus we should expect to find no persistent LMXBs above a certain critical orbital period $`P_{\mathrm{crit}}`$. For the neutron–star case this was recognised by Li & Wang (1998), who found $`P_{\mathrm{crit}}(\mathrm{NS})20`$ d, in agreement with observation. For the black–hole case, combining (3, 5) gives
$$P_{\mathrm{crit}}(\mathrm{BH})2.0f_{0.7}^{1.5}g^{0.75}\left(\frac{\dot{M}}{0.5\dot{M}_{\mathrm{Edd}}}\right)^{0.75}m_1^{1/8}m_2^{1/8}\mathrm{d},$$
(6)
where we have included a factor $`(\dot{M}/0.5\dot{M}_{\mathrm{Edd}})`$ to allow for the fact that the radiation pressure limit for the accretion rate $`\dot{M}`$ may in practice be below $`\dot{M}_{\mathrm{Edd}}`$. We thus expect to find no persistent black–hole LMXBs above this period. This is indeed supported by the available data, but hardly surprising in view of the difficulty in identifying black holes in persistent systems. Note that in high–mass black–hole systems such as Cygnus X–1, the powerful UV luminosity of the companion star, as well as the small disc size expected in a wind–fed system, are both likely to keep the disc hot and therefore give a persistent system.
## 4 GRO J1655-40
With $`g1`$ as argued above, the value of $`P_{\mathrm{crit}}(\mathrm{BH})`$ found above is close to the observed period $`P=2.62`$ d of the black–hole soft X–ray transient GRO J1655-40 (the nearest periods among alternative black–hole systems are $`P=6.47`$ d for V404 Cyg and $`P=1.23`$ d for 4U 1543–47). Indeed Kolb et al (1997) pointed out the system’s proximity to the Eddington limit during outburst, and Hynes et al. (1998) explicitly suggested that no globally steady disc solution might be possible for this system with $`\dot{M}<\dot{M}_{\mathrm{Edd}}`$. GRO J1655-40 is unusual in at least two respects:
1. The companion star has spectral type F3 – F6IV and mass $`M_22.3\mathrm{M}_{}`$. On a conventional view, this places it in the Hertzsprung gap. The companion star should therefore be expanding on a thermal timescale and thus driving a mass transfer rate $`\dot{M}_210^7\mathrm{M}_{}\mathrm{yr}^1`$ (Kolb et al., 1997). This is well above the appropriate value of $`\dot{M}_{\mathrm{crit}}^{\mathrm{irr}}`$, making it puzzling that the system is nevertheless transient, and far above the mean mass accretion rate of $`\dot{M}_{\mathrm{obs}}=1.26\times 10^{10}\mathrm{M}_{}\mathrm{yr}^1`$ deduced by van Paradijs (1996) from observation. Regös, Tout & Wickramasinghe (1998) appeal to convective overshooting to increase the main–sequence radius of stars of $`2\mathrm{M}_{}`$. The companion might then be on the main sequence rather than in the Hertzsprung gap. This implies a slower evolutionary radius expansion, bringing the predicted mass transfer rate below $`\dot{M}_{\mathrm{crit}}^{\mathrm{irr}}`$. However $`\dot{M}_2`$ is still predicted to lie uncomfortably far above $`\dot{M}_{\mathrm{obs}}`$.
2. The system was first detected in an outburst in 1994, and had probably been quiescent for at least 30 yr before that. Yet two more outbursts followed in the next two years.
The considerations given here offer explanations for both of these unusual features. First, if $`P>P_{\mathrm{crit}}(\mathrm{BH})`$, the system must be transient in some sense, regardless of the actual mass transfer rate (cf Hynes et al., 1998). It would therefore not be necessary to appeal to convective overshooting. Further, since the system is close to $`P_{\mathrm{crit}}(\mathrm{BH})`$, it is evidently accreting at a value close to the Eddington rate during its quasi–steady states (see below), making it natural that $`\dot{M}_{\mathrm{obs}}`$ is much smaller than the predicted mass transfer rate $`\dot{M}_2`$.
Second, assuming that the quantity $`g`$ has a relatively constant value close to unity, as argued above, we see from (6) that the value of $`P_{\mathrm{crit}}(\mathrm{BH})`$ is most sensitive to the filling factor $`f`$ (I consider the effect of dropping the assumption $`g`$ constant below). Thus if $`f`$ decreases, $`P_{\mathrm{crit}}(\mathrm{BH})`$ can increase above the actual orbital period, allowing irradiation to keep the disc in the high state (prolong an outburst) for as long as $`f`$ remains sufficiently small. Hence the unusual outburst behaviour of GRO J1655-40 may be explicable in terms of the time evolution of the disc size. Encouragingly there is some observational evidence (see the discussion in Orosz & Bailyn, 1997) that the grazing eclipses seen in the optical are time–dependent, just as expected if the disc size varies. Moreover Soria, Wu & Hunstead (1999) find evidence from the rotational velocities of double–peaked emission lines that the disc is at some epochs slightly larger than its tidal limit. The large resultant torques on the disc suggest that this state cannot persist and the disc must eventually shrink.
In fact we do expect $`f`$ to evolve systematically: in the early part of an outburst, the central accretion of low angular–momentum material will raise the average disc angular momentum and thus cause $`f`$ to increase, hence lowering $`P_{\mathrm{crit}}(\mathrm{BH})`$ and making the system more vulnerable to a return to quiescence. However at some stage matter transferred from the companion will tend to reduce the angular momentum of the outer disc, thus decreasing $`f`$, raising $`P_{\mathrm{crit}}(\mathrm{BH})`$ and allowing irradiation to stabilize the disc in the high state. But eventually the disc must grow towards its tidal limit, increasing $`f`$ and thus lowering $`P_{\mathrm{crit}}(\mathrm{BH})`$ again, finally enforcing a return to quiescence. Obviously a full disc code is required to follow this sequence in detail and to check if it can account qualitatively for the unusual outburst behaviour of GRO J1655–40.
Clearly, systematic evolution of one or more of the quantities appearing in $`g`$ during the outburst could have a similar effect in making $`P_{\mathrm{crit}}(\mathrm{BH})`$ oscillate around the actual orbital period $`P`$. The most likely alternative candidate is the disc aspect ratio $`H/R`$, which would appear explicitly with the power 1.5 if we substitute for $`g`$ in (6). The aspect ratio could evolve systematically on a viscous timescale because the disc may warp under radiative torques (Pringle, 1996). A warp presenting more of the disc surface to the central source would tend to stabilize it against a return to quiescence even though the central luminosity was below the Eddington limit. Again considerably more work is required to check this possibility.
## 5 Conclusions
I have shown that the Eddington limit implies a critical orbital period $`P_{\mathrm{crit}}(\mathrm{BH})`$ beyond which black–hole LMXBs cannot appear as persistent systems. The precise value of $`P_{\mathrm{crit}}(\mathrm{BH})`$ is subject to uncertainties expressed by the quantity $`g`$ in (3). I have argued that $`g`$ cannot be very far from unity if we are to understand the difference in the stability properties of discs in neutron–star and black–hole systems. In this case GRO J1655-40 lies much closer to $`P_{\mathrm{crit}}(\mathrm{BH})`$ than any other black–hole system.
The unusual behaviour of GRO J1655-40 may result from its location very close to $`P_{\mathrm{crit}}(\mathrm{BH})`$; evolution of the disc size or possible radiative warping may move the system across the boundary where a sub–Eddington luminosity can keep the disc stably in the high state. This system, and those at longer orbital periods, probably have central accretion rates which are highly super–Eddington during outbursts. Since observed radiative luminosities are mildly sub–Eddington, most of this mass must be expelled. Strong support for this comes from the observation of P Cygni profiles in GRO J1655–40 (Hynes et al., 1998). The superluminal jets observed (Hjellming & Rupen, 1995) in an outburst of this system may therefore simply represent the most dramatic part of this outflow.
## 6 Acknowledgment
I gratefully acknowledge the support of a PPARC Senior Fellowship. |
warning/0002/astro-ph0002192.html | ar5iv | text | # Strong, Variable Circular Polarization in PKS 1519-273
## 1. Introduction
Circular polarization (CP) in extragalactic sources is very small, typically 0.05% to 0.1% of the total source flux density (e.g. Roberts et al. 1975, Seaquist et al. 1974, Weiler and de Pater 1983) and sometimes variable (Komesaroff et al., 1984). Recent measurements of CP in extragalactic sources have rekindled debate as to its characteristics and origin. Wardle et al. (1998) detected CP in 3C 279 and attributed it to the presence of a relativistic pair plasma. New evidence has emerged that Sgr A, the AGN-like object at the core of our own Galaxy, is also weakly circularly polarized (Bower et al. 1999, Sault & Macquart 1999).
We present Australia Telescope Compact Array (ATCA) measurements of the timescale and magnitude of the variability of the CP in the extragalactic, intraday variable (IDV) BL Lac object, PKS 1519-273 (White et al. 1988). PKS 1519-273, at Galactic co-ordinates $`l=339.5^{},b=24.5^{}`$, is identified with a $`m_V=18.5`$ star-like object with a featureless optical spectrum. The lower limit on its redshift is $`z=0.2`$ (Veron-Cetty & Veron 1993). PKS 1519-273 is a compact high brightness temperature radio source (Linfield et al. 1989). The ATCA IDV Survey data shows strong IDV (Kedziora-Chudczer 1998) and IDV of the total and polarized flux densities at GHz frequencies has been found during each of the 5 epochs of ATCA observations over the past 7 years.
PKS 1519-273 has not been seen to exhibit IDV at either optical (Heidt & Wagner 1996) or mm wavelengths (Steppe et al. 1988, Steppe et al. 1995). However, it does have a high degree ($`512`$%) of variable optical linear polarization (Impey and Tapia 1988). PKS 1519-273 is a weak, soft X-ray source with a flux density at 1 keV of 0.39 $`\mu `$Jy (Urry et al. 1996). Its $`\gamma `$-ray energy output is less than $`0.7\times 10^7`$photons cm<sup>-2</sup>sec<sup>-1</sup> for energies $`E>100`$ MeV (Fichtel et al. 1994).
## 2. Observations and Results
We base our present report on PKS 1519-273 on the data obtained with ATCA over 5 days starting on 1998 September 9. Data were collected simultaneously for two frequencies centered on either 1.384 and 2.496 GHz, or 4.800 and 8.640 GHz each with a 128 MHz bandwidth. To ensure high quality amplitude and phase calibration we frequently observed both the standard primary flux density calibrator, PKS 1934-638, and a secondary calibrator, PKS 1514-241. The primary calibrator was used to determine accurately the flux density scale and the instrumental polarization leakages (e.g. Sault, Killeen & Kesteven 1991). The total and polarized flux density lightcurves of PKS 1519-273 are presented in fig. 1. The circularly polarized emission, is unresolved on all ATCA baselines and is strongly variable at 4.8 and 8.6 GHz. Comparison of the 2.5, 4.8 and 8.6 GHz Stokes $`V`$ measurements for PKS 1519-273 and the strong calibrator source PKS 1514-241, extensive testing and consistency checks demonstrate that the observed CP and its variations are not instrumental effects.
The most striking features of the 4.8 and 8.6 GHz light curves in fig. 1 are the exceptionally high level of CP and the large amplitude variability in all four Stokes parameters. The fractional variability of both the circularly and linearly polarized flux density exceeds that of the total flux density. The high degree of correlation between the fluctuations in $`I`$ and $`V`$ (see figs. 1 & 2) suggests that the mechanism of variability of the CP is strongly related to that in $`I`$. Comparison of the fluctuations in $`V`$ with those in $`I`$ implies that, although the overall CP is only $`1`$%, the CP of the variable component, $`\mathrm{\Delta }V/\mathrm{\Delta }I`$, is $`2.4\pm 1.3`$%, $`3.8\pm 0.4`$% and $`2.6\pm 0.5`$% at 2.4, 4.8 and 8.6 GHz respectively (see figs. 2 & 3(c)). The CP is weaker, $`<1.3`$% at 1.4 GHz and its variability is less well-established.
## 3. Discussion
### 3.1. Scintillation
We attribute the short-timescale variability of this source to Interstellar Scintillation (ISS) in our Galaxy. ISS has already been invoked to explain radio source IDV (Heeschen & Rickett 1997), including the rapid variability of PKS 0405-385 (Kedziora-Chudczer et al. 1997). If intrinsic to the source, the total intensity variations observed at 4.8 GHz imply a brightness temperature of $`T_B3\times 10^{17}`$ K for $`z0.2`$, based simply on light travel times. However, assuming that the variations are intrinsic implies a source size that is necessarily sufficiently small to exhibit variability due to ISS (Rickett et al. 1995). This suggests that an explanation based on ISS should be sought first.
The increase in modulation index (the rms normalized by the mean intensity), shown in fig. 3, and the short variability timescale from 8.6 to 4.8 GHz, shown in fig. 1, are consistent with scintillation in the regime of weak scattering (e.g., Narayan 1992), while the decrease in modulation indices and the increasingly longer variability timescales from 4.8 to 1.4 GHz are characteristic of refractive scintillation.
Assuming that the density inhomogeneities in the ISM are located on a thin screen, the refractive scintillation at 1.4 GHz may be used to place a constraint upon the distance to the scattering screen. The physical extent of the scattering disk at 1.4 GHz is the product of the long-period, refractive variability timescale, $`t_{1.4}`$, no less than 4 days (see fig. 1), and the scintillation speed, $`v`$, of order 50 km/s (see Rickett et al. 1995). VSOP observations at 1.7 GHz<sup>1</sup><sup>1</sup>1See http://www.vsop.isas.ac.jp/general/pr/1519-273.gif indicate that the source is unresolved, so we assume an angular size of no more than 0.3 mas. This implies an observer-screen distance of $`D390(v/50\mathrm{km}/\mathrm{s})(t_{1.4}/4\mathrm{days})`$ pc, or 390 pc in the present case.
Having obtained a lower limit to the distance to the scattering screen, we may constrain the angular diameter of the source from the scintillation parameters in the weak scattering regime where the scattering is quite sensitive to source size effects. The scintillation timescale of $`12`$ hours at 4.8 GHz can be explained either in terms a scattering screen at large distance ($`>15`$ kpc), or by a partially resolved source. For weak scattering, the source is resolved if the angular diameter of the source, $`\theta _S`$, exceeds the angular diameter of the first Fresnel zone $`\theta _F=(kD)^{1/2}`$, where $`k`$ is the wavenumber. The scintillation timescale is then $`t_{4.8}\theta _SD/v`$ for $`\theta _S>\theta _F`$ (Narayan 1992). A screen in our own Galaxy implies $`D15`$ kpc, so the source must be partially resolved. Assuming the asymptotic results of weak scattering to be valid between the weak and strong scattering regimes, the scintillation timescale then yields an estimate of the intrinsic angular source size of
$$\theta _S14.4\left(\frac{t_{4.8}}{12\mathrm{hours}}\right)\left(\frac{v}{50\mathrm{km}/\mathrm{s}}\right)\left(\frac{D}{1\mathrm{kpc}}\right)^1\mu \mathrm{as}.$$
$`(1)`$
For a scintillating source of flux density $`I_0`$, the root-mean-square fluctuation is $`I_{\mathrm{rms}}=I_0m(\theta _S)`$, where $`m(\theta _S)=(kD\theta _S^2)^{7/12}`$ is the modulation index expected for a source of size $`\theta _S>\theta _F`$ (e.g. Narayan 1992), and we have assumed that 4.8 GHz is near the transition frequency between weak and strong scattering. Given $`I_{\mathrm{rms}}=0.11`$ Jy and with the derived angular size of the scintillating component of the source, we estimate $`I_0`$ and derive its brightness temperature:
$$T_b2.0\times 10^{14}\left(\frac{D}{1\mathrm{kpc}}\right)^{17/12}\left(\frac{t_{4.8}}{12\mathrm{hours}}\right)^{5/6}\left(\frac{v}{50\mathrm{km}/\mathrm{s}}\right)^{5/6}\mathrm{K}.$$
$`\left(2\right)`$
Using the limit of the distance to the scattering screen, the maximum possible angular size of the source for $`t_{4.8}=12`$ hours and $`v=50`$ km/s is $`37\mu `$as, the minimum brightness temperature is $`T_b=5\times 10^{13}`$ K, consistent with incoherent synchrotron emission subject to relativistic beaming with a Doppler boosting factor $`\delta 200(1+z)`$ (Readhead 1994).
However, if the CP observed at 4.8 GHz is entirely due to the variable component, we may further constrain $`T_b`$. From fig. 2 we have $`I_0=0.35\pm 0.04`$ Jy, implying $`m(\theta _S)0.32`$ (consistent with the modulation index observed in $`V`$: $`0.32`$) and hence an angular size of $`9.8(D/1\mathrm{kpc})^{1/2}\mu `$as. Comparing with equation (1), we have $`v=34(D/1\mathrm{kpc})^{1/2}`$ km/s, implying a brightness temperature of $`3\times 10^{14}(D/1\mathrm{kpc})(t_{4.8}/12\mathrm{hours})^{5/6}`$ K. A scintillation speed exceeding $`50`$ km/s therefore implies a brightness temperature $`T_b6\times 10^{14}`$ K.
### 3.2. Circular Polarization
We consider the origin of the CP in terms of intrinsic synchrotron (Legg & Westfold 1968), partial conversion of the linear polarization into CP due to ellipticity of the natural wave modes of the cold background plasma (Pacholczyk 1973) or of the relativistic electron gas itself (e.g. Sazonov 1969, Jones & O’Dell 1977a,b).
If the CP in the scintillating component is due to synchrotron emission then, following Legg & Westfold (1968), the Lorentz factor of the particles responsible for a CP of $`m_c`$ in a uniform magnetic field is $`\gamma \mathrm{cot}\theta /m_c`$, where $`\theta `$ is the angle between the magnetic field and the line of sight. For $`30^{}<\theta <60^{}`$ this implies Lorentz factors in the range $`1545`$ to explain the CP at 4.8 GHz. Indeed, the observed (low) level of linear polarization suggests a non-uniform magnetic field, indicating that even higher Lorentz factors are required. The maximum brightness temperature of such emission is $`T_B5.9\times 10^9\gamma `$ K$`=2.7\times 10^{11}`$ K in the rest frame. Bulk motion with a Doppler boosting factor $`\delta 200`$ would account for the difference between the rest-frame and scintillation-derived brightness temperatures. Such a Doppler factor, although very high, cannot be entirely ruled out. However the observed frequency dependence of the degree of CP is far from the $`\nu ^{1/2}`$ expected from synchrotron theory (see fig 3c); the CP decreases sharply between 8.6 and 1.4 GHz. It is therefore unlikely that the CP is due to synchrotron emission.
Alternatively, the CP could be due to propagation through a relativistic pair plasma, such as may be present within the source itself. The birefringence induced in a medium by the presence of a pair plasma may convert linear polarization to CP as follows:
$$V(\nu )=U_0(\nu )\mathrm{sin}(c^3\mathrm{RRM}/\nu ^3),$$
$`(3)`$
where the relativistic rotation measure, RRM, depends upon the density of relativistic particles, the path length, the magnetic field, and the minimum Lorentz factor of the pairs (e.g. Kennett & Melrose 1998). This effect operates only when the direction of the incident linear polarization is at an oblique angle to the projection of the magnetic field on the plane orthogonal to the ray direction. The axes used to define the Stokes parameters may be chosen such that synchrotron emission has $`Q0,U=0`$. With this choice, the effect occurs only if the incident radiation has $`U_00`$, requiring either Faraday rotation or that it originate from a region of the source where the magnetic field is in a different direction to that where the polarization conversion takes place.
A characteristic of this model is a strong frequency dependence on the sign of $`V`$. If $`\mathrm{RRM}`$ is high enough to produce the CP observed at high frequency, this model predicts rapid changes in $`V`$ at low frequency: in particular, equation (3) yields the lower limit $`|\mathrm{RRM}|6.2\times 10^2(1+z)^3/f_U(8.6\mathrm{GHz})`$ rad/m<sup>3</sup> to explain the $`2.6`$% CP at 8.6 GHz, where we write $`f_U(8.6\mathrm{GHz})=|(U_0(8.6\mathrm{GHz})/I(8.6\mathrm{GHz})|`$. For this lower limit, $`\lambda ^3\mathrm{RRM}`$ will vary by $`0.88/f_U(8.6\mathrm{GHz})`$ rad across 64 MHz bandwidth at 1.4 GHz and $`0.08/f_U(8.6\mathrm{GHz})`$ rad at 2.5 GHz. We searched for frequency-dependent variations in $`V`$ at all four frequencies by selecting two adjacent 32 MHz sub-bands at each frequency. None were found. Although the Faraday rotation (RM $``$ 69 rad/m<sup>2</sup>) across the band was clearly detected at 1.4 and 2.5 GHz, the variations in $`V`$ between sub-bands were less than 4% and 0.8% at these frequencies respectively. This result appears inconsistent with the derived lower limit on RRM, although the null result at 1.4 GHz may result from an absence of CP in the scintillating component (which in turn implies $`|U_0|/I|V|/I<0.01`$ at 1.4 GHz).
The fact that all detections of the CP are of same sign also argues against this model. If (i) $`V`$ does not change sign at any frequencies intermediate to those of our measurements (i.e. the spectrum is well-sampled) and (ii) $`U_0`$ does not change sign in the range 1.4$``$8.6 GHz then equation (3) implies $`\lambda ^3\mathrm{RRM}<2\pi `$ for all frequencies above 1.4 GHz (even if $`V(1.4\mathrm{GHz})`$, whose sign is uncertain, is positive). At 1.4 GHz one then has $`|\mathrm{RRM}|<6.2\times 10^2(1+z)^3`$ rad/m<sup>3</sup>, requiring $`f_U(8.6\mathrm{GHz})100`$% to be consistent with the lower limit on $`\mathrm{RRM}`$ obtained above. While difficult to exclude entirely, we therefore conclude that the production of CP by a pair-dominated plasma is implausible.
Polarization conversion may also occur in a medium containing a mixture of cold and relativistic pair plasma, in which case the fractional CP varies as $`m_c\nu ^1`$ Pacholczyk (1973). This model appears implausible in light of the observed $`\nu ^{0.7_{0.3}^{+1.4}}`$ frequency dependence of the CP from 1.4 to 4.8 GHz.
The presence of several distinct sub-components may alter the spectral properties of the observed CP. However, the scintillation characteristics argue against the existence of multiple circularly polarized components, each with distinct $`U_0`$, RRM and $`\gamma `$. The presence of multiple components with different $`V/I`$ would lead to substructure in the lightcurve of $`V`$ compared to the lightcurve of $`I`$ as the scintillation selectively amplifies and deamplifies parts of the source differently. This would result in a loss of correlation between $`V`$ and $`I`$, particularly at 4.8 and 8.6 GHz, where the scintillation is most sensitive to small-scale structure. This is not observed in fig. 1 where the correlation coefficients are close to unity. However, it is more difficult to ascertain the presence of substructure in the variability at 1.4 and 2.5 GHz due to the long timescale of the fluctuations.
Jones & O’Dell (1977a, 1977b) presented a model for the CP of inhomogeneous synchrotron sources, incorporating optical depth effects, mode coupling and mode conversion due to the birefringence of the plasma. Below the self-absorption turnover frequency, $`\nu _{\mathrm{SSA}}`$, mode coupling dominates, and the CP is typically less than 0.05%, and certainly not more than 2%. Mode conversion dominates above $`\nu _{\mathrm{SSA}}`$, with the CP as high as 10% near $`\nu _{\mathrm{SSA}}`$, and decreasing to less than $`10^3`$ at frequencies a decade above $`\nu _{\mathrm{SSA}}`$. This model is viable only if $`\nu _{\mathrm{SSA}}`$ is within a factor $`1.4`$ of the frequency at which the high (3.8%) CP was observed, at 4.8 GHz. This is difficult to verify as we do not know the intrinsic spectrum of the scintillating component and the frequency range of our observations is limited.
## 4. Conclusion
The variability detected in PKS 1519-273 in all four Stokes parameters at frequencies from 1.4 to 8.6 GHz is remarkable, but has a natural interpretation in terms of ISS. The scintillation properties at 4.8 GHz constrain the brightness temperature of the scintillating component to $`T_b5\times 10^{13}`$ K, although there is strong evidence to suggest it may be as high as $`6\times 10^{14}`$ K. Comparison of the fluctuations in $`I`$ and $`V`$ imply that this component is exceptionally highly circularly polarized at 8.6 and 4.8 GHz. Simple applications of synchrotron theory and models of circular repolarization encounter difficulties with the spectral behavior and magnitude of the CP. The strong correlation between the fluctuations in $`I`$ and $`V`$ at 4.8 and 8.6 GHz and the high sensitivity of the scintillation to source structure at these frequencies argue against a complex source, with different $`V/I`$ in each component. Inclusion of effects due to small-scale inhomogeneity, mode coupling and optical depth effects may reproduce the observed characteristics of the CP. However, this model is only viable if the frequency at which the source is observed to become optically thin is in the range $`3.4\mathrm{GHz}\nu 6.7\mathrm{GHz}`$. This possibility is presently difficult to confirm. Even if correct, the puzzle remains as to why so few sources exhibit such high levels of CP.
Finally, in light of the extremely high brightness temperature of PKS 1519-273, we advance the possibility that the observed emission is not due to synchrotron emission at all and that high CP may be a characteristic of a new emission mechanism.
We thank Don Melrose, Ron Ekers, Mark Walker, Jim Lovell, Dick Hunstead and Lawrence Cram for valuable discussions. The Australia Telescope is funded by the Commonwealth Government for operation as a national facility by the CSIRO. |
warning/0002/hep-ex0002057.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Possible non-covariance of laws of nature was searched intensively during this century. Experimentally no evidence of covariance breaking in general relativity and electrodynamics was observed in the famous Eötvös and Hughes-Drever experiments . These precise experiments narrowed the field of possible non-covariant models, but still left a room for them.
Non-covariance, in particular Lorentz non-invariance, was proposed to be searched in weak decays by Rédei in 1966 . He extrapolated to weak decays the argument of Blokhintsev , who pointed out that possible existence of the universal length parameter violates the Lorentz invariance by introducing a preferred Lorentz frame, where this parameter is measured. In Rédei’s approach the lifetime of weakly decaying particles turns out to be anomalously dependent on their boost:
$$\tau (v)=\gamma \tau _0(1+\gamma ^2a_0^2),$$
(1)
where $`\tau _0`$ is the lifetime in the preferred frame and $`a_0`$ is the universal length parameter. Experimentally an indication of the anomalous dependence of the lifetime of charged pions and kaons on their boost was observed, while a strict limit for such a dependence was obtained for muon decays . Since the covariance violation mechanism was supposed to be the same for muons and hadrons the latter measurements were considered to rule out Rédei’s hypothesis at a high accuracy level.
In the end of 70-th the question of non-covariance in weak decays was renewed by Nielsen and Picek . They proposed not a vector but a tensor covariance breaking mechanism by introducing a non-covariant metric into the Higgs kinematic term $`h^{\mu \nu }D_\mu \varphi (D_\nu \varphi )^{},`$ where $`h^{\mu \nu }`$ is different from the metric tensor $`g^{\mu \nu }`$. Such a term results in non-covariant observables in weak decays, while electrodynamics and gravitational sectors remain covariant. Indeed, the new metric is revealed in a gauge boson mass term ($`h^{\mu \nu }A_\mu ^aA_\nu ^a\varphi ^2`$) and results in a tensor structure of the Fermi constant: $`G_F^{\mu \nu }=\frac{\sqrt{2}g_2^2}{8h^{\mu \nu }\varphi ^2}\frac{\sqrt{2}g_2^2}{8m_W^2}(g^{\mu \nu }+\xi ^{\mu \nu })`$. In the metric $`\xi ^{\mu \nu }`$ is assumed to be isotropic but not Lorentz invariant. In the simplest form it is parametrized by a single parameter $`\alpha `$ ($`\xi ^{\mu \nu }=\alpha \mathrm{diag}(1,1/3,1/3,1/3)`$). Under these assumptions the dependence of the lifetime of weakly decaying particles on their boost is derived. It turns out that the lifetime of charged pions and kaons is affected by the additional boost term in a way similar to the Rédei approach, while for muons the effect cancels out.
In the last decade the arguments in favor of Lorentz invariance violation came from the string theory, where non-local objects – strings can lead to the spontaneous breaking of the covariance . The covariance-violating term is generated when tensor rather than scalar fields gain the vacuum expectation values. If this tensor field couples to the weak gauge bosons we come to the approach of Nielsen and Picek.
The metric $`\xi ^{\mu \nu }`$ being anisotropic leads to visible anisotropy in weak decays, which can be searched for experimentally. Assuming the simplest one-parameter case ($`\chi _{\mu \nu }=\alpha \mathrm{diag}(1,0,0,1)`$) the flow rate of daughter particles in weak decays is calculated to be direction dependent. Consider for example muon decay at rest.
$`\mu ^+(p)e^+(k)+\nu _e(q_1)+\overline{\nu }_\mu (q_2).`$ (2)
Following calculations from non-covariant term in muon differential width is equal to:
$`{\displaystyle \frac{d\mathrm{\Gamma }}{dkd\mathrm{cos}\theta d\varphi }}={\displaystyle \frac{G_F^2k\chi _{\mu \nu }}{24\pi ^4m}}(2q^2p^\mu k^\nu +(pk)q^\mu q^\nu (kq)p^\mu q^\nu `$
$`(pq)k^\mu q^\nu )={\displaystyle \frac{G_F^2k\alpha }{24\pi ^4m}}(m^3(m3k)k^2m^2\mathrm{cos}(2\theta )),`$ (3)
where $`q`$ is the sum of the two neutrino momenta $`q=q_1+q_2=pk`$ and $`\theta `$ is the angle of the electron momentum with respect to $`z`$-axis (further referred as “preferred axis”). The integration over electron momentum and polar angle $`\varphi `$ gives:
$$\frac{d\mathrm{\Gamma }}{d\mathrm{cos}\theta }=(1+2\alpha \mathrm{cos}(2\theta ))\frac{d\mathrm{\Gamma }_{SM}}{d\mathrm{cos}\theta }.$$
(4)
The $`\beta `$-decays of neutron and nuclei are calculated in a similar way. The $`\beta `$-electron rate exhibits a similar directional dependence:
$$\frac{d\mathrm{\Gamma }}{d\mathrm{cos}\theta }=(1+A\alpha \mathrm{cos}(2\theta ))\frac{d\mathrm{\Gamma }_{SM}}{d\mathrm{cos}\theta },$$
(5)
where $`A`$ is $`𝒪(1)`$ and depends on nuclear form-factors.
An indication of the anisotropy of the light propagation through the Universe published in gives another argument in favor of isotropy violation, though on macro scales. In particle physics an evidence of the directional dependence of the $`\beta `$-decay rate was reported recently . In this paper we present our study of such a dependence. The upper limit for the spatial anisotropy obtained here is much stricter than the effect reported in .
## 2 Experimental Setup
The dependence of $`\beta `$-decay of $`Sr^{90}`$ on the direction of the electron emission was studied. To detect $`\beta `$-electrons two different options were used: scintillators viewed by photomultipliers and pad silicon detectors. The detectors were placed in front of a radioactive source of an intensity of $`15`$mCi corresponding to $`10^8`$ decays per second at a distance of $`4`$cm (figure 1).
The plastic scintillator had an area of 12$`\times 20`$mm<sup>2</sup> and a thickness of $`10`$mm. The Hamamtsu-5600U photomultiplier was attached to a side of the scintillator. Signal from the photomultiplier exceeded the discriminator threshold of $`30`$mV (corresponding to electron energy $`200`$keV) was counted by commercial CAMAC scaler during some exposure time. The timing was provided by a quartz timing unit with a time stability of better than $`10^5`$.
The silicon detector consisted of 4 rectangular pads with an area size of $`15\times 20`$mm<sup>2</sup> and a thickness of $`300\mu `$. Signals from the silicon pads were amplified by operational amplifiers. Under high electron rates with moderate time characteristics of the amplifier, we had to integrate the signals, rather than count them. Thus the $`\beta `$-electrons energy deposition was measured with silicon detector. Signals were integrated by RC-circuit with integration time of $`250\mu `$sec and read out with a frequency of $`5`$kHz by commercial ADC. The successive readings thus were not totally independent, their correlations were taking into account while the data was analyzed.
In order to measure the decay rate in different directions, it was possible to use the daily rotation of the Earth, keeping the measuring device at a fixed position (as in ). However, in this case a very high stability of the measurement conditions should be provided during days. High voltage and threshold potentials could change because of daily temperature and humidity variations resulting in a variation of the decay rate count with time which is a source of irreducible systematic error. The required stability seems unrealistic for the measurements with the accuracy of the order of $`10^5`$. Therefore the decision was to rotate the experimental setup artificially with minimal period such that the environment parameters could not change significantly. The Monte Carlo simulation with the actual characteristics of our device shows that the rotation with a period of a few minutes guaranteed the systematic changes of the count rate to be much smaller than the statistical error of its measurement.
The shielded source, detector, preamplifiers and power supply units were installed on the platform, rotated by an electric motor in the horizontal plane. The number of counts from the PM or the integrated charge from the silicon detector were measured during exposure time of $`4`$sec, then the platform was rotated by $`30^0`$ and the measurement was repeated. The number of steps of single $`30^0`$ rotations in one direction was equal to 10 scanning the angle of $`300^0`$, then the measurements were repeated rotating in the opposite direction. The measurements were carried on continuously during 9 days.
## 3 Data analysis and results
With the first option of the electron detector (photomultiplier) we faced the irreducible source of systematic error. Since the amplification of the PM is influenced by the external magnetic field and the vector of the Earth magnetic field changed relative to the PM while rotated, we observed a non-uniform behavior of the count rate at different directions. To reduce this effect we used active compensation of the Earth magnetic field which allowed to suppress it by a factor of $`100`$. The measurements showed that even with these special efforts, the observed nonuniformity of the PM counts was of the order of $`10^4`$. The count rate dependence on the position of the rotating platform is shown in figure 2 for two hours of data
taking. The well seen sinusoidal behavior is explained by the dependence of PM amplification on magnetic field projection. This was proved by changing the vector of the external magnetic field using the active magnetic compensation. Thus this option was used only for checking of our sensitivity to the non-uniform effects.
The silicon detectors are not affected by the magnetic field at the required level of accuracy. From the other hand the silicon detector and preamplifiers are much more sensitive to the temperature variations as demonstrated by figure 3. The integrated charge in the silicon detector depending on time during 9 days of data taking is presented in figure 3. The prominent periodical local maxima and minima
are due to day-night temperature variation of about $`3^0`$C. This instability resulted in only small error in non-uniformity measurements using the frequent rotations for different direction scanning.
The statistical error of the measurement of the integrated charge during one step of measurement was calculated from the statistical fluctuations of two subsequent steps of measurements. The differences of values obtained in all pairs of successive measurements were plotted and than fitted with Gaussian function. For cross check we also derive the statistical error from the RMS of the ADC readings during one step of measurement, taking into account the correlation between two successive readings of charge from ADC (reading with interval of $`200\mu `$sec while the integration time of RC-circuit is $`250\mu `$sec). Both ways gave the similar values within $`5\%`$. Finally we confirm the extracted from the data values of the statistical errors by numerical calculations.
Each day of data taking was divided into 12 time intervals of $`2`$hours. During each interval the Earth orientations was assumed to be the same. For each time interval the charge collected on four silicon pads for each orientation was summed up taking into account consecutive shifts in the position of each pad. Thus, the analyzed data contains information from 10 points of different orientation of the device relative to the Earth and 12 points of different Earth orientations. The data is presented in the form of 12 histograms in figure 4. Each of 12 histogram contains the dependence of the collected charge on the angle of the rotation of the device for some Earth position. All distributions was normalized to unity and unity was then subtracted, thus demonstrating only the net studied effect. Nowhere a signal of non-uniformity is seen.
The orientation dependent behavior of the $`\beta `$-electron energy deposition assuming the model discussed in the introduction section is the following:
$$\frac{\mathrm{d}N}{\mathrm{d}\mathrm{cos}\theta }1+A(t)\mathrm{cos}(2\theta +\varphi _0(t)),$$
(6)
where the angle $`\theta `$ is the angle of the rotation of our device with respect to the axis south-north. The amplitude $`A(t)`$ and the phase of the cosine $`\varphi _0(t)`$ change with time because of the Earth rotation as shown in figure 5.
Different lines represent $`A(t)`$ and $`\varphi _0(t)`$ behavior for different angles between the axis of the Earth rotation and the preferred axis for the latitude of the place of the experiment. One concludes that whatever the orientation of the preferred axis in the Universe, the non-uniformity of the electron flow is visible at least sometimes during the twenty four hours.
To set an upper limit for the dependence of electron energy on direction, these 12 histograms were fit simultaneously by a function with three free parameters: amplitude of the effect and two variables to describe direction of the preferred vector in the Universe. For each of 12 histograms the fitting function corresponds to the formula (6), while $`A(t)`$ and $`\varphi (t)`$ are different for different histograms and are functions of three parameters described above.
The fit gives the value $`A=(6.7\pm 3.6)10^6`$ for the amplitude of the effect. The upper limit, derived from likelihood function, is calculated to be $`1.410^5`$ at the $`90\%`$ confidence level.
## 4 Summary
Dependence of the rate and energy deposition of $`\beta `$-electrons from $`Sr^{90}`$ decays on the direction of emission was investigated. Unlike the previously stated evidence for such a dependence in , no signal of the non-uniformity in $`\beta `$-electron flow was observed. The upper limit on the amplitude of non-uniform behavior of $`1.410^5`$ was obtained.
Acknowledgments
It is our great pleasure to thank M.Danilov for support, useful discussions and interest to the work. We are very grateful to K.Boreskov, L.Okun and M.Voloshin for numerous theoretical discussions. We thank I.Tikhomirov, L.Laptin and A.Petryaev for their help in preparation of the experiment. |
warning/0002/hep-th0002071.html | ar5iv | text | # References
1. Introduction
The superparticle due to Siegel (later referred to as $`AB`$ superparticle) has originally been proposed as a way to avoid problematic second class constraints intrinsic to a conventional superparticle (superstring) theory without sacrificing manifest Lorentz covariance. To compensate a mismatch in the number of degrees of freedom between the $`AB`$ model and the conventional formulation , it was suggested to introduce further ($`C`$ and $`D`$) constraints, the total set forming a closed algebra. The equivalence of the $`ABCD`$ theory to the conventional superparticle (superstring) has been claimed in Ref. , thus suggesting an intriguing alternative to the standard formalism<sup>3</sup><sup>3</sup>3It has to be mentioned, however, that only in the string case the proof seems to be completely consistent. For the mechanics analogue the higher order fermionic $`C`$ constraints were treated differently (a la Gupta–Bleuler) from others.. Yet, although in the modified theories constraints do form a closed algebra and are straightforward to be realized quantum mechanically (operator quantization) , the results of the path integral quantization seem to be intractable because of infinite reducibility of the fermionic constraints involved.
A common way to attack the latter problem, which proved to be successful for the original superparticle and superstring (field theory applications were proposed earlier in Ref. ), is to make use of Lorentz harmonics to extract linearly independent components from the fermionic constraints in a covariant fashion. For the $`AB`$ model this has been accomplished in Ref. yielding a theory of rank two after a proper $`BRST`$ treatment. It has to be noted, however, that the approach of Ref. is essentially Hamiltonian. Moreover, the standard spin–statistics relations do not hold for some of the variables involved. Another serious problem is the noncompactness of the coset space parametrized by the harmonics used (see the discussion in Ref. ).
Recently, an alternative technique to cure the infinite ghost tower problem intrinsic to the Siegel superparticle, superstring has been proposed in Ref. . The idea was to appropriately extend the original phase space and then effectively cancel the infinite ghost tower by that coming from the sector of auxiliary variables.
In the present paper we investigate in full details this, looking somewhat exotic, possibility and show that the result of the quantization correlates well with that obtained previously within the framework of the harmonic superspace approach . The advantage of the novel scheme, however, is the presence of an explicit Lagrangian formulation and the standard spin–statistic relations which hold for all the variables. For simplicity of the presentation we restrict ourselves to the four dimensional case. With some modifications, however, this can be generalized to other dimensions .
In the next section we review the $`AB`$ model in $`d=4`$. It is shown that, in contrast to the common opinion, a ghost free and unitary quantum mechanics can be constructed if one sacrifices conventional conjugation properties for fermionic operators and chooses a specific modification. A light–cone Hilbert space is explicitly constructed which can be identified with the one particle sector of the quantized supersymmetric massless Wess–Zumino model. In Sec. 3, following the ideology of an earlier work , we embed the original Siegel model into an appropriately extended configuration space. In the Hamiltonian framework, the variables from the auxiliary sector turn out to be subjected to reducible constraints like those entering Siegel’s theory. These can further be used to put the complete constraint set into the form of first stage of reducibility, the latter admitting straightforward path integral quantization . Sec. 4 contains the construction of the $`BRST`$ charge in the minimal ghost sector. The theory in the extended phase space proves to be rank two. This is in perfect agreement with the analysis of the alternative harmonic superspace approach. An extension to the nonminimal ghost sector and the path integral quantization are accomplished in Sec. 5. We conclude with some remarks on possible further developments of the formalism in Sec. 6. Appendix contains the light–cone notation and technical points related to the Hamiltonian analysis of Sec. 3. Throughout the paper we use the spinor notation from Ref. .
2. Review of the $`4d`$ $`AB`$ model. No negative norm states and unitarity in the physical subspace.
Retaining only first class constraints of the conventional superparticle<sup>4</sup><sup>4</sup>4Conventions adopted in this section are $`(\theta ^\alpha )^{}=\overline{\theta }^{\dot{\alpha }}`$, $`(p_{\theta \alpha })^{}=p_{\overline{\theta }\dot{\alpha }}`$. (see Ref. for the details of the Dirac procedure)
$$p^2=0,(p_\theta \sigma ^np_n)_{\dot{\alpha }}=0,(\sigma ^np_{\overline{\theta }}p_n)_\alpha =0,$$
(1)
where $`(p_n,p_{\theta \alpha },p_{\overline{\theta }\dot{\alpha }})`$ are momenta conjugate to the configuration space variables $`(x^n,\theta ^\alpha ,\overline{\theta }^{\dot{\alpha }})`$, one discovers the constraint set to describe Siegel’s superparticle . The corresponding canonical Hamiltonian reads
$$H=\frac{1}{2}p^2.$$
(2)
Owing to the null vector $`p_n`$ entering the problem, only half of the fermionic constraints is linearly independent. In particular, the identity
$$(p_\theta \sigma ^np_n)_{\dot{\alpha }}Z_{1}^{}{}_{}{}^{\dot{\alpha }\alpha }+Z_{1}^{}{}_{}{}^{\alpha }p^20,$$
(3)
where $`Z_{1}^{}{}_{}{}^{\dot{\alpha }\alpha }=(\stackrel{~}{\sigma }^np_n)^{\dot{\alpha }\alpha },Z_{1}^{}{}_{}{}^{\beta }=p_{\theta }^{}{}_{}{}^{\beta }`$, is satisfied. On the constraint surface not all of the functions $`Z_{1}^{}{}_{}{}^{\dot{\alpha }\alpha }`$ prove to be independent
$$Z_{1}^{}{}_{}{}^{\dot{\alpha }\alpha }Z_{2}^{}{}_{\alpha \dot{\beta }}{}^{}0,Z_{2}^{}{}_{\alpha \dot{\beta }}{}^{}=(\sigma ^np_n)_{\alpha \dot{\beta }}.$$
(4)
Apparently, this process can be continued, the system at hand being infinite stage of reducibility following the terminology of Ref. .
Proceeding to the light–cone analysis of the model, one imposes the conventional gauge in the fermionic sector ($`A^\pm =\pm \frac{1}{\sqrt{2}}(A^0\pm A^3)`$)
$$\theta \sigma ^+=0,\sigma ^+\overline{\theta }=0,$$
(5)
or, reducing this to components
$$\theta ^2=0,\overline{\theta }^{\dot{2}}=0.$$
(6)
The partially reduced phase space includes then three pairs<sup>5</sup><sup>5</sup>5In what follows we omit the indices carried by the Fermi variables. $`(x^n,p_n)`$,$`(\theta ,p_\theta )`$,$`(\overline{\theta },p_{\overline{\theta }})`$, these obeying usual canonical commutation relations and the conjugation properties ($`x^n`$, $`p_n`$ are real)
$$(\theta )^{}=\overline{\theta },(p_\theta )^{}=p_{\overline{\theta }}.$$
(7)
The gauge fixed Hamiltonian action acquires the form
$$S=𝑑\tau \{p_m\dot{x}^m+p_\theta \dot{\theta }+p_{\overline{\theta }}\dot{\overline{\theta }}\frac{1}{2}p^2\}.$$
(8)
Going over to the quantum description $`(\{\widehat{\theta },\widehat{p_\theta }\}=i,\{\widehat{\overline{\theta }},\widehat{p_{\overline{\theta }}}\}=i)`$, it is customary to require a scalar product in a Hilbert space to respect the conjugation property (7), i.e.
$$\widehat{\theta }^+=\widehat{\overline{\theta }},\widehat{p_\theta }^+=\widehat{p_{\overline{\theta }}}.$$
(9)
This, however, immediately leads us to the conclusion that there are negative norm states in such a quantum space. Actually, introducing the operators
$`\widehat{a}=\frac{1}{\sqrt{2}}(\widehat{\theta }i\widehat{p}_{\overline{\theta }}),\widehat{a}^+=\frac{1}{\sqrt{2}}(\widehat{\overline{\theta }}i\widehat{p}_\theta ),`$
$`\{\widehat{a},\widehat{a}^+\}=1,`$
$`\widehat{b}=\frac{1}{\sqrt{2}}(\widehat{\theta }+i\widehat{p}_{\overline{\theta }}),\widehat{b}^+=\frac{1}{\sqrt{2}}(\widehat{\overline{\theta }}+i\widehat{p}_\theta ),`$
$`\{\widehat{b},\widehat{b}^+\}=1,`$ (10)
with a representation space being a tensor product of the corresponding Fock spaces, one discovers a ghost state due to the last line in Eq. (S0.Ex4).
It does not seem to have been emphasized previously, that a ghost free quantum mechanics still can be constructed if we sacrifice Eq. (9) and choose the alternative
$$\widehat{p}_{\theta }^{}{}_{}{}^{+}=i\widehat{\theta },\widehat{p}_{\overline{\theta }}^{}{}_{}{}^{+}=i\widehat{\overline{\theta }}.$$
(11)
With such a choice, the operators
$`\widehat{a}=\frac{1}{\sqrt{2}}(\widehat{\theta }i\widehat{p}_{\overline{\theta }}),\widehat{a}^+=\frac{1}{\sqrt{2}}(\widehat{\overline{\theta }}i\widehat{p}_\theta ),`$
$`\widehat{b}=\frac{1}{\sqrt{2}}(\widehat{\theta }+i\widehat{p}_{\overline{\theta }}),\widehat{b}^+=\frac{1}{\sqrt{2}}(\widehat{\overline{\theta }}+i\widehat{p}_\theta ),`$ (12)
obey
$$\{\widehat{a},\widehat{a}^+\}=1,\{\widehat{b},\widehat{b}^+\}=1,$$
(13)
and the corresponding Fock space, obviously, does not involve ghosts. It is worth mentioning that, there is no any physical obstruction to define a conjugation as in Eq. (11) because eigenvalues of the Fermi operators are odd supernumbers and do not correspond to any physical observables. Notice also that the gauge fixed action (8) remains to be real under the modified conjugation (11), provided the integration by parts has been performed (one can easily check that the variation problem is not influenced by the conjugation since the fermions satisfy the first order free equations).
An explicit representation of the operators $`(\widehat{\theta },\widehat{\overline{\theta }},\widehat{p}_\theta ,\widehat{p}_{\overline{\theta }})`$ in a quantum space with a scalar product respecting Eq. (11) has been given in Ref. (similar issues have been discussed in Ref. ). This is realized on a linear span of the four vectors $`(|0,|,|,|)`$, which we collectively call $`|\sigma `$, with $`(\widehat{\theta },\widehat{\overline{\theta }},\widehat{p}_\theta ,\widehat{p}_{\overline{\theta }})`$ operating like
$`\widehat{\theta }|0=0,\widehat{\theta }|=i|0,\widehat{\theta }|=0,\widehat{\theta }|=i|,`$
$`\widehat{\overline{\theta }}|0=0,\widehat{\overline{\theta }}|=0,\widehat{\overline{\theta }}|=i|0,\widehat{\overline{\theta }}|=i|,`$
$`\widehat{p_\theta }|0=|,\widehat{p_\theta }|=0,\widehat{p_\theta }|=|,\widehat{p_\theta }|=0,`$
$`\widehat{p_{\overline{\theta }}}|0=|,\widehat{p_{\overline{\theta }}}|=|,\widehat{p_{\overline{\theta }}}|=0,\widehat{p_{\overline{\theta }}}|=0,`$ (14)
and
$$\sigma |\sigma ^{}=\delta _{\sigma \sigma ^{}}.$$
(15)
The total Hilbert space is defined to be a tensor product of the linear span and the space of square integrable functions on which $`\widehat{x}^n`$ and $`\widehat{p}_n`$ act in the usual coordinate representation.
A physical Hilbert space in the complete quantum space is specified by the only constraint remaining
$$\widehat{p}^2|\mathrm{phys}=0.$$
(16)
Restricting ourselves to momentum eigenfunction, one finds
$`\mathrm{\Phi }_{p,\sigma }=\frac{1}{\sqrt{2(2\pi )^3}}|\sigma e^{ip^0t+i\stackrel{}{p}\stackrel{}{x}},`$ (17)
where for physical reasons we have chosen an upper shell of the light cone $`p^0=\sqrt{\stackrel{}{p}^2}`$.
A scalar product in the physical subspace is given by
$$\mathrm{\Phi }|\mathrm{\Psi }=id^3\stackrel{}{x}(\overline{\mathrm{\Phi }}_0\mathrm{\Psi }_0\overline{\mathrm{\Phi }}\mathrm{\Psi }),$$
(18)
or
$$\mathrm{\Phi }_{p,\sigma }|\mathrm{\Phi }_{p^{},\sigma ^{}}=p^0\delta ^{(3)}(\stackrel{}{p}\stackrel{}{p}^{^{}})\delta _{\sigma \sigma ^{}},$$
(19)
for the momentum eigenfunctions.
It is instructive then to clarify the structure of the Pauli-Lubanski vector for the case at hand. Putting the classical expression
$$W_a=\frac{1}{2}ϵ_{abcd}p^bS^{cd},$$
(20)
with $`S^{cd}=\theta ^\delta \left(\sigma ^{cd}\right)_{\delta }^{}{}_{}{}^{\gamma }p_{\theta \gamma }+p_{\overline{\theta }\dot{\gamma }}\left(\stackrel{~}{\sigma }^{cd}\right)_{}^{\dot{\gamma }}{}_{\dot{\delta }}{}^{}\overline{\theta }^{\dot{\delta }}`$ being the spin part of the Lorentz generators, onto the surface of the constraints and gauges, one obtains
$$W_a=\frac{i}{2}p_a\left(p_\theta \theta p_{\overline{\theta }}\overline{\theta }\right).$$
(21)
Here we made use of the identities ($`ϵ_{0123}=1`$)
$$\sigma _{ab}=\frac{i}{2}ϵ_{abcd}\sigma ^{cd},\stackrel{~}{\sigma }_{ab}=\frac{i}{2}ϵ_{abcd}\stackrel{~}{\sigma }^{cd}.$$
(22)
Owing to the minus sign between the two terms entering Eq. (21), one does not face any operator ordering ambiguities in passing to quantum description. In particular,
$$\widehat{W}_a\mathrm{\Phi }_{p,\sigma }=\sigma p_a\mathrm{\Phi }_{p,\sigma },$$
(23)
where the number coefficient $`\sigma `$ takes values
$$\sigma =(0,\frac{1}{2},\frac{1}{2},0),$$
(24)
for the states $`|\sigma =(|0,|,|,|)`$, respectively. Observe also that $`\widehat{W}_a`$ is hermitian with respect to both conjugation prescriptions (9),(11).
Since the construction of unitary irreducible representations (irreps) of the Poincaré group reduces to that of the little group generated by $`\widehat{W}_a`$ (see e.g. ), it is straightforward to verify that given the vector $`|\sigma =(|0,|,|,|)`$ in Eq. (17) the corresponding linear space ($`p_a`$ takes values on the upper shell of the light cone) realizes a unitary irrep of helicity $`\sigma `$, with $`\sigma `$ being specified in Eq. (24).
Finally, it is worth mentioning, that the set of helicities obtained allows us to identify the quantum space constructed with the one particle sector of a quantized supersymmetric massless Wess–Zumino model. This correlates well with the results of the Dirac quantization<sup>6</sup><sup>6</sup>6In Ref. quantum wave functions were realized on real scalar superfields. This condition can be weakened to include complex scalar superfields if one makes proper use of both of the equation entering Eq. (34b) of Ref. . accomplished in Ref. .
A path integral representation for the superpropagator of the massless Wess–Zumino model that explicitly involves a gauge fixed action of the $`4d`$ Siegel superparticle has been given in Ref. .
3. Siegel superparticle in an extended phase space
As was demonstrated in the previous section, the Siegel superparticle in the original formulation is infinite stage of reducibility. In this section we reformulate the model by introducing a set of auxiliary variables. The extension makes it possible to put fermionic constraints into an irreducible form valid for subsequent path integral quantization.
3.1. Action and symmetries
The superparticle action to be analyzed is
$`S={\displaystyle }d\tau \{{\displaystyle \frac{1}{2e}}(\dot{x}^m+i\theta \sigma ^m\dot{\overline{\theta }}i\dot{\theta }\sigma ^m\overline{\theta }+i\psi \sigma ^m\overline{\rho }i\rho \sigma ^m\overline{\psi }+\omega \mathrm{\Lambda }^m)^2`$
$`\rho ^\alpha \dot{\theta }_\alpha \overline{\rho }_{\dot{\alpha }}\dot{\overline{\theta }}^{\dot{\alpha }}\omega \varphi \mathrm{\Lambda }^2\mathrm{\Lambda }_mi\phi \sigma ^m\overline{\chi }+\mathrm{\Lambda }_mi\chi \sigma ^m\overline{\phi }\}.`$ (25)
The theory is invariant under the standard rigid supersymmetry transformations. The local reparametrizations and $`\kappa `$-symmetry of Siegel’s model
$$\begin{array}{ccc}\delta _\alpha \theta =\alpha \dot{\theta },\hfill & \delta _\alpha \overline{\theta }=\alpha \dot{\overline{\theta }},\hfill & \delta _\alpha x^n=\alpha \dot{x}^n,\hfill \\ \delta _\alpha \rho =\alpha \dot{\rho },\hfill & \delta _\alpha \overline{\rho }=\alpha \dot{\overline{\rho }},\hfill & \delta _\alpha e=(\alpha e)^{},\hfill \\ \delta _\alpha \psi =(\alpha \psi )^{},\hfill & \delta _\alpha \overline{\psi }=(\alpha \overline{\psi })^{},\hfill & \delta _\alpha \omega =(\alpha \omega )^{},\hfill \\ \delta _\alpha \mathrm{\Lambda }^n=\alpha \dot{\mathrm{\Lambda }}^n,\hfill & \delta _\alpha \chi =\alpha \dot{\chi },\hfill & \delta _\alpha \overline{\chi }=\alpha \dot{\overline{\chi }},\hfill \\ \delta _\alpha \phi =(\alpha \phi )^{},\hfill & \delta _\alpha \overline{\phi }=(\alpha \overline{\phi })^{},\hfill & \delta _\alpha \varphi =(\alpha \varphi )^{},\hfill \end{array}$$
$$\begin{array}{c}\delta _\kappa \theta =ie^1\mathrm{\Pi }_n\sigma ^n\overline{\kappa },\delta _\kappa \overline{\theta }=ie^1\mathrm{\Pi }_n\kappa \sigma ^n,\hfill \\ \delta _\kappa x^n=i\delta _\kappa \theta \sigma ^n\overline{\theta }i\theta \sigma ^n\delta _\kappa \overline{\theta }i\kappa \sigma ^n\overline{\rho }+i\rho \sigma ^n\overline{\kappa },\hfill \\ \delta _\kappa e=4\dot{\theta }\kappa +4\overline{\kappa }\dot{\overline{\theta }},\delta _\kappa \psi =\dot{\kappa },\hfill \\ \delta _\kappa \overline{\psi }=\dot{\overline{\kappa }},\hfill \end{array}$$
$`(26)`$
where $`\mathrm{\Pi }^m=\dot{x}^m+i\theta \sigma ^m\dot{\overline{\theta }}i\dot{\theta }\sigma ^m\overline{\theta }+i\psi \sigma ^m\overline{\rho }i\rho \sigma ^m\overline{\psi }+\omega \mathrm{\Lambda }^m`$, are extended by two new symmetries depending on fermionic parameters $`\beta ,\gamma `$, the latter acting in the sector of the new variables
$$\delta _\beta \chi =\overline{\beta }\stackrel{~}{\sigma }^n\mathrm{\Lambda }_n,\delta _\beta \overline{\chi }=\mathrm{\Lambda }_n\stackrel{~}{\sigma }^n\beta ,\delta _\beta \varphi =i(\phi \beta \overline{\phi }\overline{\beta }),$$
(27)
$$\delta _\gamma \phi =\overline{\gamma }\stackrel{~}{\sigma }^n\mathrm{\Lambda }_n,\delta _\gamma \overline{\phi }=\mathrm{\Lambda }_n\stackrel{~}{\sigma }^n\gamma ,\delta _\gamma \varphi =i(\chi \gamma \overline{\chi }\overline{\gamma }).$$
(28)
From the transformation rules above, one concludes that the variables $`(x^m,\theta ^\alpha ,\overline{\theta }_{\dot{\alpha }})`$ parametrize a conventional $`R^{4|4}`$ superspace, $`(e,\psi ^\alpha ,\overline{\psi }_{\dot{\alpha }})`$ prove to be gauge fields for local reparametrizations and $`\kappa `$–symmetry, whereas the pair $`(\rho ^\alpha ,\overline{\rho }_{\dot{\alpha }})`$ provides the terms corresponding to a (mixed) covariant propagator for fermions. This holds as in the Siegel model. As shown below, there is no dynamics in the sector of the new variables $`(\omega ,\mathrm{\Lambda }^m,\varphi ,\phi ^\alpha ,\overline{\phi }_{\dot{\alpha }},\chi ^\alpha ,\overline{\chi }_{\dot{\alpha }})`$, these prove to be purely auxiliary.
3.2. Fermionic constraints made irreducible
Proceeding to the Hamiltonian analysis one finds fourteen primary constraints<sup>7</sup><sup>7</sup>7We define momenta conjugate to Fermi variables to be right derivatives of a Lagrangian with respect to velocities. This corresponds to the following choice of the Poisson brackets $`\{\theta ^\alpha ,p_{\theta }^{}{}_{\beta }{}^{}\}=\delta _{}^{\alpha }{}_{\beta }{}^{},\{\overline{\theta }_{\dot{\alpha }},p_{\overline{\theta }}^{}{}_{}{}^{\dot{\beta }}\}=\delta _{\dot{\alpha }}^{}{}_{}{}^{\dot{\beta }}`$ and the position of momenta and velocities in the Hamiltonian as specified below in Eq. (S0.Ex13). Our conventions for the conjugation slightly differ from those used in the review section $`(\theta ^\alpha )^{}=\overline{\theta }^{\dot{\alpha }}`$, $`(p_{\theta \alpha })^{}=p_{\overline{\theta }\dot{\alpha }}`$.
$`p_e=0,p_\psi =0,p_{\overline{\psi }}=0,p_\rho =0,p_{\overline{\rho }}=0,p_\omega =0,`$
$`p_\mathrm{\Lambda }=0,p_\varphi =0,p_\phi =0,p_{\overline{\phi }}=0,p_\chi =0,p_{\overline{\chi }}=0,`$
$`p_{\theta \alpha }p_ni(\sigma ^n\overline{\theta })_\alpha \rho _\alpha =0,p_{\overline{\theta }}^{}{}_{}{}^{\dot{\alpha }}+p_ni(\theta \sigma ^n)^{\dot{\alpha }}\overline{\rho }^{\dot{\alpha }}=0,`$ (29)
where $`p_q`$ stands for a momentum canonically conjugate to a variable $`q`$. The total Hamiltonian has the form
$`H=p_e\lambda _e+p_{\psi \alpha }\lambda _{\psi }^{}{}_{}{}^{\alpha }+p_{\overline{\psi }}^{}{}_{}{}^{\dot{\alpha }}\lambda _{\overline{\psi }\dot{\alpha }}+p_{\rho \alpha }\lambda _{\rho }^{}{}_{}{}^{\alpha }+p_{\overline{\rho }}^{}{}_{}{}^{\dot{\alpha }}\lambda _{\overline{\rho }\dot{\alpha }}+p_\omega \lambda _\omega +p_{\mathrm{\Lambda }n}\lambda _{\mathrm{\Lambda }}^{}{}_{}{}^{n}+p_\varphi \lambda _\varphi `$
$`+p_{\phi \alpha }\lambda _{\phi }^{}{}_{}{}^{\alpha }+p_{\overline{\phi }}^{}{}_{}{}^{\dot{\alpha }}\lambda _{\overline{\phi }\dot{\alpha }}+p_{\chi \alpha }\lambda _{\chi }^{}{}_{}{}^{\alpha }+p_{\overline{\chi }}^{}{}_{}{}^{\dot{\alpha }}\lambda _{\overline{\chi }\dot{\alpha }}+(p_{\overline{\theta }}+p_ni\theta \sigma ^n\overline{\rho })^{\dot{\alpha }}\lambda _{\overline{\theta }\dot{\alpha }}`$
$`+(p_\theta p_ni\sigma ^n\overline{\theta }\rho )_\alpha \lambda _{\theta }^{}{}_{}{}^{\alpha }+\frac{1}{2}ep^2i\psi \sigma ^n\overline{\rho }p_n+i\rho \sigma ^n\overline{\psi }p_n+\varphi \mathrm{\Lambda }^2+\omega (1\mathrm{\Lambda }p)`$
$`+i\phi \sigma ^n\overline{\chi }\mathrm{\Lambda }_ni\chi \sigma ^n\overline{\phi }\mathrm{\Lambda }_n,`$ (30)
where $`\lambda _{\mathrm{}}`$ are the Lagrange multipliers associated to the primary constraints. The conservation in time of the primary constraints yields the secondary ones
$`p^2=0,p_n(\sigma ^n\overline{\rho })_\alpha =0,p_n(\rho \sigma ^n)_{\dot{\alpha }}=0,`$
$`\mathrm{\Lambda }_n(\sigma ^n\overline{\chi })_\alpha =0,\mathrm{\Lambda }_n(\chi \sigma ^n)_{\dot{\alpha }}=0,`$
$`\mathrm{\Lambda }_n(\sigma ^n\overline{\phi })_\alpha =0,\mathrm{\Lambda }_n(\phi \sigma ^n)_{\dot{\alpha }}=0,`$
$`\mathrm{\Lambda }^2=0,1\mathrm{\Lambda }p=0,`$
$`2\varphi \mathrm{\Lambda }^n+\omega p^ni\phi \sigma ^n\overline{\chi }+i\chi \sigma ^n\overline{\phi }=0,`$ (31)
and fixes some of the Lagrange multipliers,
$$\begin{array}{cc}\lambda _\theta =p_ni\sigma ^n\overline{\psi },\hfill & \lambda _{\overline{\theta }}=i\psi \sigma ^np_n,\hfill \\ \lambda _\rho =2p_ni\sigma ^n\lambda _{\overline{\theta }}0,\hfill & \lambda _{\overline{\rho }}=2i\lambda _\theta \sigma ^np_n0.\hfill \end{array}$$
(32)
Beautifully enough, the last equation in Eq. (S0.Ex16) can be simplified to (a proof is given in Appendix)
$$\omega =0,2\varphi i\phi \sigma ^n\overline{\chi }p_n+i\chi \sigma ^n\overline{\phi }p_n=0.$$
(33)
With this remark, consistency conditions for the secondary constraints amount to
$`p\lambda _\mathrm{\Lambda }=0,\mathrm{\Lambda }\lambda _\mathrm{\Lambda }=0,\lambda _\omega =0,`$
$`2\lambda _\varphi =i\lambda _\phi \sigma ^n\overline{\chi }p_n+i\lambda _\chi \sigma ^n\overline{\phi }p_ni\phi \sigma ^n\lambda _{\overline{\chi }}p_n+i\chi \sigma ^n\lambda _{\overline{\phi }}p_n,`$
$`\mathrm{\Lambda }_n(\sigma ^n\lambda _{\overline{\chi }})_\alpha +\lambda _{\mathrm{\Lambda }n}(\sigma ^n\overline{\chi })_\alpha =0,\mathrm{\Lambda }_n(\lambda _\chi \sigma ^n)_{\dot{\alpha }}+\lambda _{\mathrm{\Lambda }n}(\chi \sigma ^n)_{\dot{\alpha }}=0,`$
$`\mathrm{\Lambda }_n(\sigma ^n\lambda _{\overline{\phi }})_\alpha +\lambda _{\mathrm{\Lambda }n}(\sigma ^n\overline{\phi })_\alpha =0,\mathrm{\Lambda }_n(\lambda _\phi \sigma ^n)_{\dot{\alpha }}+\lambda _{\mathrm{\Lambda }n}(\phi \sigma ^n)_{\dot{\alpha }}=0.`$ (34)
Making use of the light–cone arguments like those given in the Appendix one can show that each of the fermionic equations entering Eq. (S0.Ex20) determines precisely half of the corresponding fermionic Lagrange multipliers.
Thus no tertiary constraints arise at this stage, the complete constraint system being
$$p_e=0,p_\psi =0,p_{\overline{\psi }}=0,$$
(35)
$$p_\rho =0,p_\theta p_ni(\sigma ^n\overline{\theta })\rho =0,$$
(36)
$$p_{\overline{\rho }}=0,p_{\overline{\theta }}+p_ni(\theta \sigma ^n)\overline{\rho }=0,$$
(37)
$$p_\omega =0,\omega =0,$$
(38)
$$p_\varphi =0,2\varphi i\phi \sigma ^n\overline{\chi }p_n+i\chi \sigma ^n\overline{\phi }p_n=0,$$
(39)
$$p_\phi =0,\phi \sigma ^n\mathrm{\Lambda }_n=0,$$
(40)
$$p_{\overline{\phi }}=0,\sigma ^n\overline{\phi }\mathrm{\Lambda }_n=0,$$
(41)
$$p_\chi =0,\chi \sigma ^n\mathrm{\Lambda }_n=0,$$
(42)
$$p_{\overline{\chi }}=0,\sigma ^n\overline{\chi }\mathrm{\Lambda }_n=0,$$
(43)
$$p^2=0,p_\theta \sigma ^np_n=0,\sigma ^np_{\overline{\theta }}p_n=0,$$
(44)
$$p_\mathrm{\Lambda }=0,\mathrm{\Lambda }^2=0,1\mathrm{\Lambda }p=0.$$
(45)
The constraints (35) are first–class. Imposing the gauge
$$e=1,\psi =0,\overline{\psi }=0,$$
(46)
which yields
$$\lambda _e=0,\lambda _\psi =0,\lambda _{\overline{\psi }}=0,$$
(47)
one can disregard the canonical pairs $`(e,p_e)`$, $`(\psi ,p_\psi )`$, $`(\overline{\psi },p_{\overline{\psi }})`$. In the same manner, the variables $`(\rho ,p_\rho )`$,$`(\overline{\rho },p_{\overline{\rho }})`$,$`(\omega ,p_\omega )`$,$`(\varphi ,p_\varphi )`$ can be omitted after introducing the Dirac bracket associated with the second class constraints (36)–(39). The Dirac brackets for the remaining variables prove to coincide with the Poisson ones.
One has to be more inventive when imposing a gauge in the sector (40), ((41)). Passing to the light–cone coordinates (see Appendix) one concludes that, due to $`\mathrm{\Lambda }^2=0`$, there is only one linearly independent component entering the last of the spinor constraints (40) ((41)), this proves to be second class, whereas the corresponding momenta include one first and one second class constraints. Beautifully enough, on account of the last of the equations (45) these can be put into covariant (redundant) form
$`p_\phi =0\{\begin{array}{cc}p_\phi \sigma ^n\mathrm{\Lambda }_n=0\hfill & \text{first class}\hfill \\ p_\phi \sigma ^np_n=0\hfill & \text{second class}\hfill \end{array}`$ (50)
Fixing a gauge is now obvious (again in a covariant and redundant form)
$$\sigma ^n\phi p_n=0,$$
(51)
which yields
$$\phi =0,$$
(52)
when combined with Eq. (40). The conservation in time of the gauge (51) yields
$$\lambda _\phi \sigma ^np_n=0.$$
(53)
Together with Eq. (S0.Ex20) this completely specifies $`\lambda _\phi `$. Note also that consistency ($`(\phi )^{}=\overline{\phi }`$) requires us to impose the complex conjugate equation
$$p_n\sigma ^n\overline{\phi }=0\overline{\phi }=0.$$
(54)
One finally concludes that there is no dynamics in the sector $`(\phi ,p_\phi )`$, $`(\overline{\phi },p_{\overline{\phi }})`$.
The same arguments apply to the variables $`(\chi ,p_\chi )`$, $`(\overline{\chi },p_{\overline{\chi }})`$. For our purposes, however, it is convenient not to impose a gauge in this sector but rather use these purely auxiliary variables to supplement Siegel’s constraints (44) up to irreducible ones. Actually, it is straightforward to check that the system (see also Ref. )
$$\overline{\mathrm{\Phi }}_{\dot{\alpha }}(p_\theta \sigma ^np_n+p_\chi \sigma ^n\mathrm{\Lambda }_n)_{\dot{\alpha }}=0,\mathrm{\Phi }_\alpha (p_n\sigma ^np_{\overline{\theta }}+\mathrm{\Lambda }_n\sigma ^np_{\overline{\chi }})_\alpha =0\text{first class},$$
(55)
$$\overline{\mathrm{\Psi }}_{\dot{\alpha }}(\chi \sigma ^n\mathrm{\Lambda }_n+p_\chi \sigma ^np_n)_{\dot{\alpha }}=0,\mathrm{\Psi }_\alpha (\mathrm{\Lambda }_n\sigma ^n\overline{\chi }+p_n\sigma ^np_{\overline{\chi }})_\alpha =0\text{second class},$$
(56)
$$p^2=0\text{first class},$$
(57)
is completely equivalent to the initial equations (42)–(44). Here the identities
$`p_{\chi }^{}{}_{}{}^{\alpha }=\frac{1}{2\mathrm{\Lambda }p}\overline{\mathrm{\Phi }}_{\dot{\alpha }}(\stackrel{~}{\sigma }^mp_m)^{\dot{\alpha }\alpha }\frac{1}{2\mathrm{\Lambda }p}\overline{\mathrm{\Psi }}_{\dot{\alpha }}(\stackrel{~}{\sigma }^m\mathrm{\Lambda }_m)^{\dot{\alpha }\alpha }\frac{1}{2\mathrm{\Lambda }p}p^2p_{\theta }^{}{}_{}{}^{\alpha }\frac{1}{2\mathrm{\Lambda }p}\mathrm{\Lambda }^2\chi ^\alpha ,`$ (58)
$`p_{\overline{\chi }}^{}{}_{}{}^{\dot{\alpha }}=\frac{1}{2\mathrm{\Lambda }p}(\stackrel{~}{\sigma }^mp_m)^{\dot{\alpha }\alpha }\mathrm{\Phi }_\alpha \frac{1}{2\mathrm{\Lambda }p}(\stackrel{~}{\sigma }^m\mathrm{\Lambda }_m)^{\dot{\alpha }\alpha }\mathrm{\Psi }_\alpha \frac{1}{2\mathrm{\Lambda }p}p^2p_{\overline{\theta }}^{}{}_{}{}^{\dot{\alpha }}\frac{1}{2\mathrm{\Lambda }p}\mathrm{\Lambda }^2\overline{\chi }^{\dot{\alpha }},`$ (59)
prove to be useful. The equivalence just stated implies also that the constraint set above is irreducible, otherwise we would have less than $`8+1`$ equations and Eqs. (55)–(57) would not be equivalent to (42)–(44) ($`8+1`$ linearly independent components).
It remains to discuss the bosonic constraints (45). Constructing a (weak) projector to the directions orthogonal to the vectors $`p^n,\mathrm{\Lambda }^n`$
$$\pi _{m}^{}{}_{}{}^{n}=\delta _{m}^{}{}_{}{}^{n}p_m\mathrm{\Lambda }^n\mathrm{\Lambda }_mp^n,$$
(60)
one can easily extract first class constraints contained in $`p_\mathrm{\Lambda }`$, the complete constraint set being
$$\stackrel{~}{p}_{\mathrm{\Lambda }m}(\pi p_\mathrm{\Lambda })_m=p_{\mathrm{\Lambda }m}(p_\mathrm{\Lambda }\mathrm{\Lambda })p_m(p_\mathrm{\Lambda }p)\mathrm{\Lambda }_m=0\text{first class},$$
(61)
$$p_\mathrm{\Lambda }p=0,\mathrm{\Lambda }^2=0,p_\mathrm{\Lambda }\mathrm{\Lambda }=0,1\mathrm{\Lambda }p=0\text{second class}.$$
(62)
In view of the identities<sup>8</sup><sup>8</sup>8Here and in what follows the symbol $``$ means an equality up to a linear combination of second class constraints.
$$\stackrel{~}{p}_\mathrm{\Lambda }\mathrm{\Lambda }0,\stackrel{~}{p}_\mathrm{\Lambda }p0,$$
(63)
one concludes that there are only two linearly independent components entering Eq. (61), the total number of constraints being sufficient to suppress dynamics in the sector. In order to explicitly decouple the first class constraints above from the fermionic second class ones (56), it suffices to redefine them like
$$\stackrel{~}{p}_{\mathrm{\Lambda }}^{}{}_{}{}^{n}=0\stackrel{~}{p}_{\mathrm{\Lambda }}^{}{}_{}{}^{n}\frac{1}{2}\chi \sigma ^n\stackrel{~}{\sigma }^mp_\chi p_m\frac{1}{2}p_{\overline{\chi }}\stackrel{~}{\sigma }^m\sigma ^n\overline{\chi }p_m=0.$$
(64)
As the Dirac bracket associated with the second class constraints is introduced, this seems to be inessential here.
It is worth mentioning, that the dynamical equivalence of the model (24) and the Siegel superparticle can be easily established if one imposes the non covariant gauge
$$\mathrm{\Lambda }^i=0,\text{i=1,2}.$$
(65)
To summarize, in the extended phase space the infinite reducibility of the constraints (44) characterizing the Siegel model can be compensated by that coming from the sector of additional variables to put the fermionic constraints into an irreducible form. Residual reducibility proves to fall in the bosonic constraints (61),(62). Being the first stage of reducibility, these admit consistent path integral quantization.
Quantization of the constraint system (55)–(57),(61),(62) will be our main concern in the next sections.
3.3. The Dirac bracket
In the presence of second class constraints both the nilpotency equation to determine the $`BRST`$ charge and that to fix the unitarizing Hamiltonian should be solved under the Dirac bracket associated with the full set of second class constraints . To construct the latter, it suffices to convert the matrix of Poisson brackets of second class constraints<sup>9</sup><sup>9</sup>9The construction proves to be more involved when second class constraints in a question are (infinitely) reducible. A recipe has been given in Ref. . Denoting the constraints collectively by $`\mathrm{\Theta }_i=(p_\mathrm{\Lambda }p,\mathrm{\Lambda }^2,p_\mathrm{\Lambda }\mathrm{\Lambda },1\mathrm{\Lambda }p,\overline{\mathrm{\Psi }}_{\dot{\alpha }},\mathrm{\Psi }_\alpha )`$ and $`\mathrm{\Gamma }_{ij}\{\mathrm{\Theta }_i,\mathrm{\Theta }_j\}`$, one finds this to be
$$\mathrm{\Gamma }_{ij}=\left(\begin{array}{cccccc}0& 2\mathrm{\Lambda }p& pp_\mathrm{\Lambda }& p^2& (\chi \sigma ^np_n)_{\dot{\beta }}& (p_n\sigma ^n\overline{\chi })_\beta \\ & & & & & \\ 2\mathrm{\Lambda }p& 0& 2\mathrm{\Lambda }^2& 0& 0& 0\\ & & & & & \\ pp_\mathrm{\Lambda }& 2\mathrm{\Lambda }^2& 0& \mathrm{\Lambda }p& (\chi \sigma ^n\mathrm{\Lambda }_n)_{\dot{\beta }}& (\mathrm{\Lambda }_n\sigma ^n\overline{\chi })_\beta \\ & & & & & \\ p^2& 0& \mathrm{\Lambda }p& 0& 0& 0\\ & & & & & \\ (\chi \sigma ^np_n)_{\dot{\alpha }}& 0& (\chi \sigma ^n\mathrm{\Lambda }_n)_{\dot{\alpha }}& 0& 4(\stackrel{~}{\sigma }^{nm})_{\dot{\alpha }\dot{\beta }}\mathrm{\Lambda }_np_m& 0\\ & & & & & \\ (p_n\sigma ^n\overline{\chi })_\alpha & 0& (\mathrm{\Lambda }_n\sigma ^n\overline{\chi })_\alpha & 0& 0& 4(\sigma ^{nm})_{\alpha \beta }\mathrm{\Lambda }_np_m\end{array}\right).$$
The corresponding superdeterminant amounts to a simple number coefficient
$$sdet\mathrm{\Gamma }_{ij}=\frac{1}{4},$$
(66)
this to be used when constructing the path integral measure in Sec. 5.
Given a supermatrix $`F=F_B+F_S`$, where $`F_B`$ and $`F_S`$ are the body and the soul respectively , the inverse supermatrix is constructed according to the rule
$$F^1=F_{B}^{}{}_{}{}^{1}+\underset{k=1}{\overset{\mathrm{}}{}}(1)^k(F_{B}^{}{}_{}{}^{1}F_S)^kF_{B}^{}{}_{}{}^{1}.$$
(67)
In our case only the first two terms entering the power series above prove to be non vanishing, the corresponding inverse supermatrix being
$$\mathrm{\Gamma }^{ij}=\frac{1}{\mathrm{\Delta }}\left(\begin{array}{cccccc}0& \mathrm{\Lambda }p& 0& 2\mathrm{\Lambda }^2& 0& 0\\ & & & & & \\ \mathrm{\Lambda }p& 0& p^2& pp_{\mathrm{\Lambda }}^{}{}_{}{}^{^{}}& \frac{1}{2}(\chi \sigma ^np_n)^{\dot{\beta }}& \frac{1}{2}(p_n\sigma ^n\overline{\chi })^\beta \\ & & & & & \\ 0& p^2& 0& 2\mathrm{\Lambda }p& 0& 0\\ & & & & & \\ 2\mathrm{\Lambda }^2& pp_{\mathrm{\Lambda }}^{}{}_{}{}^{^{}}& 2\mathrm{\Lambda }p& 0& (\chi \sigma ^n\mathrm{\Lambda }_n)^{\dot{\beta }}& (\mathrm{\Lambda }_n\sigma ^n\overline{\chi })^\beta \\ & & & & & \\ 0& \frac{1}{2}(\chi \sigma ^np_n)^{\dot{\alpha }}& 0& (\chi \sigma ^n\mathrm{\Lambda }_n)^{\dot{\alpha }}& 2(\stackrel{~}{\sigma }^{nm})^{\dot{\alpha }\dot{\beta }}\mathrm{\Lambda }_np_m& 0\\ & & & & & \\ 0& \frac{1}{2}(p_n\sigma ^n\overline{\chi })^\alpha & 0& (\mathrm{\Lambda }_n\sigma ^n\overline{\chi })^\alpha & 0& 2(\sigma ^{nm})^{\alpha \beta }\mathrm{\Lambda }_np_m\end{array}\right),$$
where $`\mathrm{\Delta }2((\mathrm{\Lambda }p)^2\mathrm{\Lambda }^2p^2)`$ and $`pp_{\mathrm{\Lambda }}^{}{}_{}{}^{^{}}pp_\mathrm{\Lambda }+\frac{1}{2}(\chi ^2+\overline{\chi }^2)`$.
With the $`\mathrm{\Gamma }^{ij}`$ at hand, the Dirac bracket is straightforward to build
$`\{A,B\}_D=\{A,B\}\{A,\mathrm{\Theta }_i\}\mathrm{\Gamma }^{ij}\{\mathrm{\Theta }_j,B\}.`$ (68)
Being rather involved in the general form, the bracket considerably simplifies when evaluated in specific coordinate sectors<sup>10</sup><sup>10</sup>10In what follows we omit the label $`D`$ attached to the Dirac brackets. (only the brackets to be used below are explicitly given here)
$`\{\chi ^\alpha ,p_{\chi \beta }\}=\frac{1}{2}\delta _{}^{\alpha }{}_{\beta }{}^{}\frac{2}{\mathrm{\Delta }}\mathrm{\Lambda }p(\sigma _{nm})_{\beta }^{}{}_{}{}^{\alpha }\mathrm{\Lambda }^np^m,\{\chi ^\alpha ,\chi ^\beta \}=\frac{2}{\mathrm{\Delta }}p^2(\sigma _{nm})^{\alpha \beta }\mathrm{\Lambda }^np^m,`$
$`\{p_{\chi \alpha },p_{\chi \beta }\}=\frac{2}{\mathrm{\Delta }}\mathrm{\Lambda }^2(\sigma _{nm})_{\alpha \beta }\mathrm{\Lambda }^np^m;`$
$`\{\mathrm{\Lambda }^n,p_{\mathrm{\Lambda }m}\}=\delta _{}^{n}{}_{m}{}^{}\frac{2}{\mathrm{\Delta }}\mathrm{\Lambda }p(p^n\mathrm{\Lambda }_m+\mathrm{\Lambda }^np_m)+\frac{2}{\mathrm{\Delta }}p^2\mathrm{\Lambda }^n\mathrm{\Lambda }_m+\frac{2}{\mathrm{\Delta }}\mathrm{\Lambda }^2p^np_m,`$
$`\{\mathrm{\Lambda }^n,\mathrm{\Lambda }^m\}=0,\{p_{\mathrm{\Lambda }n},p_{\mathrm{\Lambda }m}\}=\frac{2}{\mathrm{\Delta }}p^2(\mathrm{\Lambda }_np_{\mathrm{\Lambda }m}\mathrm{\Lambda }_mp_{\mathrm{\Lambda }n})+\frac{2}{\mathrm{\Delta }}pp_\mathrm{\Lambda }(p_n\mathrm{\Lambda }_m`$
$`p_m\mathrm{\Lambda }_n)+\frac{2}{\mathrm{\Delta }}p\mathrm{\Lambda }(p_{\mathrm{\Lambda }n}p_mp_{\mathrm{\Lambda }m}p_n)\frac{i}{\mathrm{\Delta }}(\chi ^2\overline{\chi }^2)ϵ_{nmkl}\mathrm{\Lambda }^kp^l;`$
$`\{\theta ^\alpha ,p_{\theta \beta }\}=\delta _{}^{\alpha }{}_{\beta }{}^{},\{\theta ^\alpha ,\theta ^\beta \}=0,\{p_{\theta \alpha },p_{\theta \beta }\}=0,\{p_n,p_m\}=0.`$ (69)
Analogously, for the cross sectors one finds (in what follows we will not need the explicit form of the brackets involving the $`x^n`$–variable, these are omitted here)
$`\{p_{\mathrm{\Lambda }n},\chi ^\alpha \}=\frac{1}{\mathrm{\Delta }}p^2\mathrm{\Lambda }_n\chi ^\alpha +\frac{1}{\mathrm{\Delta }}p^2(\chi \sigma _n\stackrel{~}{\sigma }^k\mathrm{\Lambda }_k)^\alpha +\frac{1}{\mathrm{\Delta }}p_n(\chi \sigma ^k\mathrm{\Lambda }_k\stackrel{~}{\sigma }^mp_m)^\alpha `$ (70)
$`\frac{1}{\mathrm{\Delta }}\mathrm{\Lambda }p(\chi \sigma _n\stackrel{~}{\sigma }^kp_k)^\alpha ,`$
$`\{p_{\mathrm{\Lambda }n},p_{\chi \alpha }\}=\frac{1}{\mathrm{\Delta }}\mathrm{\Lambda }^2p_n\chi _\alpha +\frac{1}{\mathrm{\Delta }}\mathrm{\Lambda }^2(\chi \sigma _n\stackrel{~}{\sigma }^kp_k)_\alpha +\frac{1}{\mathrm{\Delta }}\mathrm{\Lambda }_n(\chi \sigma ^kp_k\stackrel{~}{\sigma }^m\mathrm{\Lambda }_m)_\alpha `$
$`\frac{1}{\mathrm{\Delta }}\mathrm{\Lambda }p(\chi \sigma _n\stackrel{~}{\sigma }^k\mathrm{\Lambda }_k)_\alpha .`$
Similar relations hold for complex conjugates. In particular, to derive brackets involving $`(\overline{\chi },p_{\overline{\chi }})`$ it suffices to raise $`\alpha `$, low $`\beta `$ and then exchange them with dotted indices.
In obtaining Eqs. (S0.Ex34),(70) the following identities
$`Tr(\sigma _{ab}\sigma _{cd})=\frac{1}{2}(\eta _{ac}\eta _{bd}\eta _{ad}\eta _{bc})+\frac{i}{2}ϵ_{abcd}`$ (71)
$`Tr(\stackrel{~}{\sigma }_{ab}\stackrel{~}{\sigma }_{cd})=\frac{1}{2}(\eta _{ac}\eta _{bd}\eta _{ad}\eta _{bc})\frac{i}{2}ϵ_{abcd},`$ (72)
with $`ϵ_{0123}=1`$ and $`\eta _{nm}=diag(,+,+,+)`$, prove to be useful.
Finally, it is worth noting that, as long as the path integral quantization is concerned, the presence of the $`\delta `$–function of second class constraints in the path integral measure allows one to solve the equations on the $`BRST`$–charge and the unitarizing Hamiltonian modulo second class constraints . In particular, this will provide further simplifications in Eqs. (S0.Ex34)–(70).
3.4. The algebra of first class constraints
Having evaluated the Dirac bracket, we are now in a position to specify the algebra of the first class constraints (55), (57), (61), the corresponding structure functions to be used when constructing the $`BRST`$ charge.
Taking into account Eq. (58), the identity
$`\chi ^2=\frac{1}{\mathrm{\Lambda }p}p^2(p_\chi \chi )\frac{1}{\mathrm{\Lambda }p}\overline{\mathrm{\Psi }}_{\dot{\alpha }}(\stackrel{~}{\sigma }^mp_m\chi )^{\dot{\alpha }},`$ (73)
and the fact that according to the general recipe it suffices to know the algebra modulo second class constraints, one finds the only nontrivial brackets to be
$`\{\stackrel{~}{p}_{\mathrm{\Lambda }n},\stackrel{~}{p}_{\mathrm{\Lambda }m}\}U_{nm}^{}{}_{}{}^{k}\stackrel{~}{p}_{\mathrm{\Lambda }k}+U_{nm}p^2,\{\stackrel{~}{p}_{\mathrm{\Lambda }n},\mathrm{\Phi }_\alpha \}U_{n\alpha }^{}{}_{}{}^{\beta }\mathrm{\Phi }_\beta +U_{n\alpha }p^2,`$ (74)
$`\{\stackrel{~}{p}_{\mathrm{\Lambda }n},\overline{\mathrm{\Phi }}_{\dot{\alpha }}\}U_{n\dot{\alpha }}^{}{}_{}{}^{\dot{\beta }}\overline{\mathrm{\Phi }}_{\dot{\beta }}+U_{n\dot{\alpha }}p^2.`$ (75)
The explicit form of the structure functions involved is ($`U_{n\dot{\alpha }}^{}{}_{}{}^{\dot{\beta }}`$, $`U_{n\dot{\alpha }}`$ are obtained by complex conjugation)
$`U_{nm}^{}{}_{}{}^{k}=\frac{2}{\mathrm{\Delta }}((\mathrm{\Lambda }_np^2p_n)\delta _{m}^{}{}_{}{}^{k}(\mathrm{\Lambda }_mp^2p_m)\delta _{n}^{}{}_{}{}^{k})`$ (76)
$`U_{nm}=\frac{i}{\mathrm{\Delta }}(p_\chi \chi p_{\overline{\chi }}\overline{\chi })ϵ_{nmkl}\mathrm{\Lambda }^kp^l,`$ (77)
$`U_{n\alpha }^{}{}_{}{}^{\beta }=\frac{1}{2}(\sigma _n\stackrel{~}{\sigma }^kp_k)_{\alpha }^{}{}_{}{}^{\beta }+\frac{1}{\mathrm{\Delta }}\mathrm{\Lambda }_np^2\delta _{\alpha }^{}{}_{}{}^{\beta }+\frac{1}{\mathrm{\Delta }}(\mathrm{\Lambda }_np^2p_n)(\mathrm{\Lambda }^k\sigma _k\stackrel{~}{\sigma }^lp_l)_{\alpha }^{}{}_{}{}^{\beta },`$ (78)
$`U_{n\alpha }=\frac{1}{2}(\sigma _np_{\overline{\theta }})_\alpha \frac{1}{\mathrm{\Delta }}\mathrm{\Lambda }_n(p^k\sigma _kp_{\overline{\theta }})_\alpha +\frac{1}{\mathrm{\Delta }}(\mathrm{\Lambda }_np^2p_n)(\mathrm{\Lambda }^k\sigma _kp_{\overline{\theta }})_\alpha .`$ (79)
Worth noting also are the important algebraic properties of the functions obtained (the same holds for complex conjugates)
$`U_{nm}\mathrm{\Lambda }^m=0,U_{n\alpha }^{}{}_{}{}^{\beta }\mathrm{\Lambda }^n0,U_{n\alpha }\mathrm{\Lambda }^n0,`$
$`U_{nm}p^m=0,U_{n\alpha }^{}{}_{}{}^{\beta }p^n0,U_{n\alpha }p^n0.`$ (80)
These will be of frequent use when establishing the nilpotency of the BRST charge in the next section.
4. The BRST charge and the unitarizing Hamiltonian in the minimal ghost sector
Proceeding to the $`BRST`$ quantization, one associates a couple of canonically conjugate ghost variables to each of the first class constraints (55), (57), (61) $`(C^{\dot{\alpha }},\overline{𝒫}_{\dot{\alpha }})`$, $`(C^\alpha ,\overline{𝒫}_\alpha )`$, $`(C,\overline{𝒫})`$,$`(C^n,\overline{𝒫}_n)`$. The statistics and the ghost number are specified by the conventional prescriptions
$`ϵ(C^A)=ϵ(\overline{𝒫}^A)=ϵ_A+1,`$
$`gh(C^A)=gh(\overline{𝒫}^A)=1.`$ (81)
To compensate the overcounting in the sector $`(C^n,\overline{𝒫}_n)`$ (only two components entering the bosonic constraint (61) are linearly independent) one further introduces the secondary ghosts $`(C^1,\overline{𝒫}^1)`$, $`(C^2,\overline{𝒫}^2)`$ which obey
$`ϵ(C^{1,2})=ϵ(\overline{𝒫}^{1,2})=0,`$
$`gh(C^{1,2})=gh(\overline{𝒫}^{1,2})=2.`$ (82)
Together with the previously introduced variables these exhaust the minimal ghost sector for the model under consideration.
The $`BRST`$ charge is defined to be a solution of the nilpotency equation
$$\{\mathrm{\Omega }_{min},\mathrm{\Omega }_{min}\}0,$$
(83)
satisfying the boundary condition
$$\mathrm{\Omega }_{min}=\mathrm{\Phi }_\alpha C^\alpha +\overline{\mathrm{\Phi }}_{\dot{\alpha }}C^{\dot{\alpha }}+\stackrel{~}{p}_{\mathrm{\Lambda }n}C^n+p^2C+\overline{𝒫}_n\mathrm{\Lambda }^nC^1+\overline{𝒫}_np^nC^2+\mathrm{}.$$
(84)
The first four terms entering Eq. (84) are typical for the $`BRST`$ quantization of irreducible gauge theories. Through Eq. (83) they automatically generate the gauge algebra (74). The two remaining terms are designed to generate the identities (63) and are specific to the treatment of reducible theories.
Calculating the contribution of the boundary terms into Eq. (83)
$`\{\mathrm{\Omega }_{min},\mathrm{\Omega }_{min}\}2\overline{𝒫}_m\{\mathrm{\Lambda }^m,\stackrel{~}{p}_{\mathrm{\Lambda }n}\}C^1C^n2(U_{n\alpha }^{}{}_{}{}^{\beta }\mathrm{\Phi }_\beta +U_{n\alpha }p^2)C^\alpha C^n`$
$`2(U_{n\dot{\alpha }}^{}{}_{}{}^{\dot{\beta }}\overline{\mathrm{\Phi }}_{\dot{\beta }}+U_{n\dot{\alpha }}p^2)C^{\dot{\alpha }}C^n(U_{nm}^{}{}_{}{}^{k}\stackrel{~}{p}_{\mathrm{\Lambda }k}+U_{nm}p^2)C^mC^n+\mathrm{},`$ (85)
one can partially clarify the structure of the terms which were missing in Eq. (84). In particular, extending the ansatz (84) by means of three new contributions
$`\frac{1}{2}\overline{𝒫}_k\stackrel{~}{U}_{nm}^{}{}_{}{}^{k}C^mC^n+\overline{𝒫}_\alpha U_{n\beta }^{}{}_{}{}^{\alpha }C^\beta C^n+\overline{𝒫}_{\dot{\alpha }}U_{n\dot{\beta }}^{}{}_{}{}^{\dot{\alpha }}C^{\dot{\beta }}C^n,`$ (86)
where
$`\stackrel{~}{U}_{nm}^{}{}_{}{}^{k}=U_{nm}^{}{}_{}{}^{k}\frac{2}{\mathrm{\Delta }}p^k(\mathrm{\Lambda }_np_m\mathrm{\Lambda }_mp_n),`$
$`\stackrel{~}{U}_{nm}^{}{}_{}{}^{k}\mathrm{\Lambda }^m\frac{2}{\mathrm{\Delta }}\{\mathrm{\Lambda }^k,p_{\mathrm{\Lambda }n}\},\stackrel{~}{U}_{nm}^{}{}_{}{}^{k}p^m0,`$ (87)
one can get rid of the first term (which is a manifestation of the reducibility of the constraints) and those involving $`\stackrel{~}{p}_\mathrm{\Lambda },\mathrm{\Phi },\overline{\mathrm{\Phi }}`$
$`\{\mathrm{\Omega }_{min},\mathrm{\Omega }_{min}\}U_{nm}p^2C^mC^n2U_{n\alpha }p^2C^\alpha C^n2U_{n\dot{\alpha }}p^2C^{\dot{\alpha }}C^n`$
$`2\overline{𝒫}_\alpha U_{n\gamma }^{}{}_{}{}^{\alpha }U_{m\beta }^{}{}_{}{}^{\gamma }C^mC^nC^\beta 2\overline{𝒫}_{\dot{\alpha }}U_{n\dot{\gamma }}^{}{}_{}{}^{\dot{\alpha }}U_{m\dot{\beta }}^{}{}_{}{}^{\dot{\gamma }}C^mC^nC^{\dot{\beta }}+\mathrm{}.`$ (88)
In order to verify Eq. (S0.Ex52) one has to use the algebraic properties of the structure functions (S0.Ex47) and the Jacobi identities resulting from the constraint algebra
$`\stackrel{~}{U}_{[m\widehat{a}}^{}{}_{}{}^{k}\stackrel{~}{U}_{bc]}^{}{}_{}{}^{a}0,U_{[m\widehat{a}}^{}{}_{}{}^{k}U_{bc]}^{}{}_{}{}^{a}0,`$
$`\{\stackrel{~}{p}_{\mathrm{\Lambda }[n},\stackrel{~}{U}_{ab]}^{}{}_{}{}^{k}\}0,\{\stackrel{~}{p}_{\mathrm{\Lambda }[n},U_{ab]}^{}{}_{}{}^{k}\}0,`$
$`\{\stackrel{~}{p}_{\mathrm{\Lambda }n},U_{m\alpha }^{}{}_{}{}^{\beta }\}\{\stackrel{~}{p}_{\mathrm{\Lambda }m},U_{n\alpha }^{}{}_{}{}^{\beta }\}U_{nm}^{}{}_{}{}^{k}U_{k\alpha }^{}{}_{}{}^{\beta }0,`$
$`\{\stackrel{~}{p}_{\mathrm{\Lambda }n},U_{m\alpha }\}\{\stackrel{~}{p}_{\mathrm{\Lambda }m},U_{n\alpha }\}U_{nm}^{}{}_{}{}^{k}U_{k\alpha }0.`$ (89)
Similar equations hold for complex conjugates. Here the square bracket stands for a complete antisymmetrization of indices and a hat over an index means that it is not affected by the antisymmetrization.
It is instructive then to give an explicit form of the terms quadratic in the structure functions which enter Eq. (S0.Ex52)
$`U_{m\alpha }^{}{}_{}{}^{\beta }U_{n\beta }^{}{}_{}{}^{\gamma }U_{n\alpha }^{}{}_{}{}^{\beta }U_{m\beta }^{}{}_{}{}^{\gamma }=\{\frac{1}{\mathrm{\Delta }}(\mathrm{\Lambda }_np_m\mathrm{\Lambda }_mp_n)(\mathrm{\Lambda }_l\sigma ^l\stackrel{~}{\sigma }^kp_k)_{\alpha }^{}{}_{}{}^{\gamma }+\frac{1}{\mathrm{\Delta }}(\mathrm{\Lambda }_np_m\mathrm{\Lambda }_mp_n)\delta _{\alpha }^{}{}_{}{}^{\gamma }+`$
$`\frac{1}{\mathrm{\Delta }}\mathrm{\Lambda }_m(\sigma _n\stackrel{~}{\sigma }^kp_k)_{\alpha }^{}{}_{}{}^{\gamma }\frac{1}{\mathrm{\Delta }}\mathrm{\Lambda }_n(\sigma _m\stackrel{~}{\sigma }^kp_k)_{\alpha }^{}{}_{}{}^{\gamma }\frac{1}{\mathrm{\Delta }}(\mathrm{\Lambda }_mp^2p_m)(\sigma _n\stackrel{~}{\sigma }^k\mathrm{\Lambda }_k)_{\alpha }^{}{}_{}{}^{\gamma }+`$
$`\frac{1}{\mathrm{\Delta }}(\mathrm{\Lambda }_np^2p_n)(\sigma _m\stackrel{~}{\sigma }^k\mathrm{\Lambda }_k)_{\alpha }^{}{}_{}{}^{\gamma }+(\sigma _{nm})_{\alpha }^{}{}_{}{}^{\beta }\}p^2\mathrm{\Pi }_{mn\alpha }^\gamma p^2,`$
$`\mathrm{\Pi }_{mn\alpha }^{}{}_{}{}^{\gamma }\mathrm{\Lambda }^n0,\mathrm{\Pi }_{mn\alpha }^{}{}_{}{}^{\gamma }p^n0.`$ (90)
Being factor of $`p^2`$ this suggests a further amendment to the $`\mathrm{\Omega }_{min}`$
$`\overline{𝒫}U_{n\alpha }C^\alpha C^n+\overline{𝒫}U_{n\dot{\alpha }}C^{\dot{\alpha }}C^n+\frac{1}{2}\overline{𝒫}U_{nm}C^mC^n`$
$`\frac{1}{2}\overline{𝒫}\overline{𝒫}_\alpha \mathrm{\Pi }_{nm\beta }^{}{}_{}{}^{\alpha }C^mC^nC^\beta \frac{1}{2}\overline{𝒫}\overline{𝒫}_{\dot{\alpha }}\mathrm{\Pi }_{nm\dot{\beta }}^{}{}_{}{}^{\dot{\alpha }}C^mC^nC^{\dot{\beta }},`$ (91)
where $`\mathrm{\Pi }_{nm\dot{\beta }}^{}{}_{}{}^{\dot{\alpha }}`$ is the complex conjugate of $`\mathrm{\Pi }_{nm\beta }^{}{}_{}{}^{\alpha }`$.
Beautifully enough, by making use of the next portion of the Jacobi identities
$`\{\stackrel{~}{p}_{\mathrm{\Lambda }[a},U_{bc]}\}+U_{[a\widehat{d}}U_{bc]}^{}{}_{}{}^{d}0,`$
$`\mathrm{\Pi }_{mn\alpha }^{}{}_{}{}^{\gamma }\mathrm{\Phi }_\gamma +U_{m\alpha }^{}{}_{}{}^{\beta }U_{n\beta }U_{n\alpha }^{}{}_{}{}^{\beta }U_{m\beta }+\{\mathrm{\Phi }_\alpha ,U_{nm}\}0,`$
$`\{\stackrel{~}{p}_{\mathrm{\Lambda }[a},\mathrm{\Pi }_{mn]\beta }^{}{}_{}{}^{\alpha }\}+\mathrm{\Pi }_{[a\widehat{d}\beta }^{}{}_{}{}^{\alpha }U_{mn]}^{}{}_{}{}^{d}0,`$
$`\mathrm{\Pi }_{[mn\delta }^{}{}_{}{}^{\gamma }U_{a]\gamma }^{}{}_{}{}^{\alpha }U_{[a\delta }^{}{}_{}{}^{\gamma }\mathrm{\Pi }_{mn]\gamma }^{}{}_{}{}^{\alpha }0,`$ (92)
and their complex conjugates, one can verify the nilpotency of our ansatz, the $`BRST`$ charge in the minimal ghost sector being of the form
$`\mathrm{\Omega }_{min}=\mathrm{\Phi }_\alpha C^\alpha +\overline{\mathrm{\Phi }}_{\dot{\alpha }}C^{\dot{\alpha }}+\stackrel{~}{p}_{\mathrm{\Lambda }n}C^n+p^2C+\overline{𝒫}_n\mathrm{\Lambda }^nC^1+\overline{𝒫}_np^nC^2+`$
$`\frac{1}{2}\overline{𝒫}_k\stackrel{~}{U}_{nm}^{}{}_{}{}^{k}C^mC^n+\overline{𝒫}_\alpha U_{n\beta }^{}{}_{}{}^{\alpha }C^\beta C^n+\overline{𝒫}_{\dot{\alpha }}U_{n\dot{\beta }}^{}{}_{}{}^{\dot{\alpha }}C^{\dot{\beta }}C^n+`$
$`\overline{𝒫}U_{n\alpha }C^\alpha C^n+\overline{𝒫}U_{n\dot{\alpha }}C^{\dot{\alpha }}C^n+\frac{1}{2}\overline{𝒫}U_{nm}C^mC^n`$
$`\frac{1}{2}\overline{𝒫}\overline{𝒫}_\alpha \mathrm{\Pi }_{nm\beta }^{}{}_{}{}^{\alpha }C^mC^nC^\beta \frac{1}{2}\overline{𝒫}\overline{𝒫}_{\dot{\alpha }}\mathrm{\Pi }_{nm\dot{\beta }}^{}{}_{}{}^{\dot{\alpha }}C^mC^nC^{\dot{\beta }}.`$ (93)
For this to be real, one has to impose the following conjugation properties on the ghost variables
$`(C^\alpha )^{}=C^{\dot{\alpha }},(C^n)^{}=C^n,(C)^{}=C,(C^{1,2})^{}=C^{1,2},`$
$`(\overline{𝒫}_\alpha )^{}=\overline{𝒫}_{\dot{\alpha }},(\overline{𝒫}_n)^{}=\overline{𝒫}_n,(\overline{𝒫})^{}=\overline{𝒫},(\overline{𝒫}^{1,2})^{}=\overline{𝒫}^{1,2}.`$ (94)
Thus, within the framework of the $`BRST`$ quantization the modified formulation proves to be a theory of rank two. Our result here correlates well with that obtained previously in the alternative harmonic superspace approach . Worth noting also is that a naive limit of the expression obtained to the original phase space breaks manifest Lorenz covariance $`(\mathrm{\Lambda }^i=0,\mathrm{\Lambda }^{}=0,\mathrm{\Lambda }^+=\frac{1}{p^{}})`$, as it should.
Finally, we observe that the boundary condition which has to be imposed on the unitarizing Hamiltonian (a proper Hamiltonian treatment requires secondary constraints to be added to the initial Hamiltonian with the corresponding Lagrange multipliers)
$$H|_{C=\overline{𝒫}=0}=H_0=0,$$
(95)
automatically satisfies the needed equation
$$\{H,\mathrm{\Omega }_{min}\}0.$$
(96)
Hence, no ghost corrections to Eq. (95) are to be added, the latter fits to describe the unitarizing Hamiltonian for the case at hand (see also Ref. ).
5. Extension to the nonminimal ghost sector. Transition amplitude.
Having constructed $`\mathrm{\Omega }_{min}`$ and $`H`$, an extension to the nonminimal ghost sector is straightforward . The irreducible constraints $`\mathrm{\Phi }_\alpha `$, $`\overline{\mathrm{\Phi }}_{\dot{\alpha }}`$, $`p^2`$ can be treated in the usual way. One introduces three canonical pairs of new ghost variables $`(𝒫^\alpha ,\overline{C}_\alpha )`$,$`(𝒫^{\dot{\alpha }},\overline{C}_{\dot{\alpha }})`$, $`(𝒫,\overline{C})`$ along with the corresponding Lagrange multipliers $`(\lambda ^\alpha ,\pi _\alpha )`$,$`(\lambda ^{\dot{\alpha }},\pi _{\dot{\alpha }})`$,$`(\lambda ,\pi )`$ (the statistics and the ghost number of the new variables are given below in the Table 1). Associated with the reducible constraints $`\stackrel{~}{p}_\mathrm{\Lambda }`$ are the primary ghosts and Lagrange multipliers $`(𝒫^n,\overline{C}_n)`$, $`(\lambda ^n,\pi _n)`$, as well as the secondary ones $`(𝒫^1,\overline{C}^1)`$, $`(𝒫^2,\overline{C}^2)`$, $`(\lambda ^1,\pi ^1)`$, $`(\lambda ^2,\pi ^2)`$. A direct inspection of the structure of the ghost sector (with the use of the Table 1) shows the disbalance between the number of unphysical degrees of freedom and that of the ghosts introduced. This can be improved by introducing further “extra” ghosts . In our case these are exhausted by $`(𝒫_{(1)}^{}{}_{}{}^{1},\overline{C}_{(1)}^{}{}_{}{}^{1})`$, $`(𝒫_{(1)}^{}{}_{}{}^{2},\overline{C}_{(1)}^{}{}_{}{}^{2})`$, $`(\lambda _{(1)}^{}{}_{}{}^{1},\pi _{(1)}^{}{}_{}{}^{1})`$, $`(\lambda _{(1)}^{}{}_{}{}^{2},\pi _{(1)}^{}{}_{}{}^{2})`$. The statistics and the ghost number of the new variables are gathered in the following table
Table 1.1 Ghosts (nonminimal sector)
| | $`𝒫^\alpha `$ | $`\overline{C}_\alpha `$ | $`𝒫^{\dot{\alpha }}`$ | $`\overline{C}_{\dot{\alpha }}`$ | $`𝒫`$ | $`\overline{C}`$ | $`𝒫^n`$ | $`\overline{C}_n`$ | $`𝒫^1`$ | $`\overline{C}^1`$ | $`𝒫^2`$ | $`\overline{C}^2`$ |
| --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- |
| $`{\displaystyle }`$ $`ϵ`$ | 0 | 0 | 0 | 0 | 1 | 1 | 1 | 1 | 0 | 0 | 0 | 0 |
| $`{\displaystyle }`$ $`gh`$ | 1 | -1 | 1 | -1 | 1 | -1 | 1 | -1 | 2 | -2 | 2 | -2 |
Table 1.2 Lagrange multipliers
| | $`\lambda ^\alpha `$ | $`\pi _\alpha `$ | $`\lambda ^{\dot{\alpha }}`$ | $`\pi _{\dot{\alpha }}`$ | $`\lambda `$ | $`\pi `$ | $`\lambda ^n`$ | $`\pi _n`$ | $`\lambda ^1`$ | $`\pi ^1`$ | $`\lambda ^2`$ | $`\pi ^2`$ |
| --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- |
| $`{\displaystyle }`$ $`ϵ`$ | 1 | 1 | 1 | 1 | 0 | 0 | 0 | 0 | 1 | 1 | 1 | 1 |
| $`{\displaystyle }`$ $`gh`$ | 0 | 0 | 0 | 0 | 0 | 0 | 0 | 0 | 1 | -1 | 1 | -1 |
Table 1.3 Extra ghosts
| | $`𝒫_{(1)}^{}{}_{}{}^{1}`$ | $`\overline{C}_{(1)}^{}{}_{}{}^{1}`$ | $`𝒫_{(1)}^{}{}_{}{}^{2}`$ | $`\overline{C}_{(1)}^{}{}_{}{}^{2}`$ | $`\lambda _{(1)}^{}{}_{}{}^{1}`$ | $`\pi _{(1)}^{}{}_{}{}^{1}`$ | $`\lambda _{(1)}^{}{}_{}{}^{2}`$ | $`\pi _{(1)}^{}{}_{}{}^{2}`$ |
| --- | --- | --- | --- | --- | --- | --- | --- | --- |
| $`{\displaystyle }`$ $`ϵ`$ | 1 | 1 | 1 | 1 | 0 | 0 | 0 | 0 |
| $`{\displaystyle }`$ $`gh`$ | 1 | -1 | 1 | -1 | 0 | 0 | 0 | 0 |
A continuation of $`\mathrm{\Omega }_{min}`$ to the complete relativistic phase space is now easy to perform
$`\mathrm{\Omega }=\mathrm{\Omega }_{min}+\pi _\alpha 𝒫^\alpha +\pi _{\dot{\alpha }}𝒫^{\dot{\alpha }}+\pi 𝒫+\pi _n𝒫^n+\pi ^1𝒫^1+\pi ^2𝒫^2+`$
$`\pi _{(1)}^{}{}_{}{}^{1}𝒫_{(1)}^{}{}_{}{}^{1}+\pi _{(1)}^{}{}_{}{}^{2}𝒫_{(1)}^{}{}_{}{}^{2}.`$ (97)
This supplies us with the last tool needed for quantizing the theory, the corresponding transition amplitude is given by the formal path integral<sup>11</sup><sup>11</sup>11As usual, the fermionic $`\delta `$–function is defined as $`\delta (\theta )=\theta `$. Hence, $`\delta (\mathrm{\Psi }_\alpha )\mathrm{\Psi }_1\mathrm{\Psi }_2\mathrm{\Psi }^2`$.
$$Z_\mathrm{\Psi }=\frac{1}{2}D\mu \delta (\mathrm{\Lambda }^2)\delta (1\mathrm{\Lambda }p)\delta (p_\mathrm{\Lambda }p)\delta (p_\mathrm{\Lambda }\mathrm{\Lambda })\delta (\mathrm{\Psi }_\alpha )\delta (\overline{\mathrm{\Psi }}_{\dot{\alpha }})e^{\frac{i}{h}S}.$$
(98)
Here the effective quantum action has the form
$`S={\displaystyle }d\tau (p_n\dot{x}^n+p_{\theta \alpha }\dot{\theta }^\alpha +p_{\overline{\theta }}^{}{}_{}{}^{\dot{\alpha }}\dot{\overline{\theta }}_{\dot{\alpha }}+p_{\mathrm{\Lambda }n}\dot{\mathrm{\Lambda }}^n+p_{\chi \alpha }\dot{\chi }^\alpha +p_{\overline{\chi }}^{}{}_{}{}^{\dot{\alpha }}\dot{\overline{\chi }}_{\dot{\alpha }}+\pi _\alpha \dot{\lambda }^\alpha +\pi ^{\dot{\alpha }}\dot{\lambda }_{\dot{\alpha }}+`$
$`\pi \dot{\lambda }+\pi _n\dot{\lambda }^n+\pi ^1\dot{\lambda }^1+\pi ^2\dot{\lambda }^2+\pi _{(1)}^{}{}_{}{}^{1}\dot{\lambda }_{(1)}^{}{}_{}{}^{1}+\pi _{(1)}^{}{}_{}{}^{2}\dot{\lambda }_{(1)}^{}{}_{}{}^{2}+\overline{𝒫}_\alpha \dot{C}^\alpha +\overline{𝒫}_{\dot{\alpha }}\dot{C}^{\dot{\alpha }}+`$
$`\overline{𝒫}\dot{C}+\overline{𝒫}_n\dot{C}^n+\overline{𝒫}^1\dot{C}^1+\overline{𝒫}^2\dot{C}^2+\overline{C}_\alpha \dot{𝒫}^\alpha +\overline{C}_{\dot{\alpha }}\dot{𝒫}^{\dot{\alpha }}+\overline{C}\dot{𝒫}+\overline{C}_n\dot{𝒫}^n+`$
$`\overline{C}^1\dot{𝒫}^1+\overline{C}^2\dot{𝒫}^2+\overline{C}_{(1)}^{}{}_{}{}^{1}\dot{𝒫}_{(1)}^{}{}_{}{}^{1}+\overline{C}_{(1)}^{}{}_{}{}^{2}\dot{𝒫}_{(1)}^{}{}_{}{}^{2}\{\mathrm{\Psi },\mathrm{\Omega }\}_D),`$ (99)
with $`D\mu `$ being the usual Liouville measure over the full phase space and $`\mathrm{\Psi }`$ denoting the gauge fixing fermion $`(ϵ(\mathrm{\Psi })=1,gh(\mathrm{\Psi })=1)`$. Given a specific form of the latter, a number of ghost (and Lagrange multiplier) integrations can be performed explicitly. For example, we can integrate out the whole bunch of variables $`(C^{1,2},\overline{𝒫}^{1,2})`$, $`(𝒫^{1,2},\overline{C}^{1,2})`$, $`(\lambda ^{1,2},\pi ^{1,2})`$, $`(𝒫_{(1)}^{}{}_{}{}^{1,2},\overline{C}_{(1)}^{}{}_{}{}^{1,2})`$, $`(\lambda _{(1)}^{}{}_{}{}^{1,2},\pi _{(1)}^{}{}_{}{}^{1,2})`$ by taking the following representation for $`\mathrm{\Psi }`$
$$\mathrm{\Psi }=\lambda ^1\overline{𝒫}^1+\lambda ^2\overline{𝒫}^2+\frac{1}{ϵ}C_np^n\overline{C}^1+\frac{1}{ϵ}C_n\mathrm{\Lambda }^n\overline{C}^2+\frac{1}{ϵ}\lambda _{(1)}^{}{}_{}{}^{1}\overline{C}_{(1)}^{}{}_{}{}^{1}+\frac{1}{ϵ}\lambda _{(1)}^{}{}_{}{}^{2}\overline{C}_{(1)}^{}{}_{}{}^{2}+\mathrm{\Psi }^{^{}}.$$
(100)
Here $`\mathrm{\Psi }^{^{}}`$ does not depend on the set above and $`ϵ`$ is a constant. An explicit integration which appeals to a passage to a discrete lattice attaches then four new factors to the path integral measure
$`Z_\mathrm{\Psi }=\frac{1}{2}{\displaystyle D\mu ^{^{}}\delta (\mathrm{\Lambda }^2)\delta (1\mathrm{\Lambda }p)\delta (p_\mathrm{\Lambda }p)\delta (p_\mathrm{\Lambda }\mathrm{\Lambda })\delta (\mathrm{\Psi }_\alpha )\delta (\overline{\mathrm{\Psi }}_{\dot{\alpha }})}`$ (101)
$`\delta (C_np^n)\delta (C_n\mathrm{\Lambda }^n)\delta (𝒫_np^n)\delta (𝒫_n\mathrm{\Lambda }^n)e^{\frac{i}{h}S},`$
and reduces the effective action to the relatively simple form
$`S={\displaystyle }d\tau (p_n\dot{x}^n+p_{\theta \alpha }\dot{\theta }^\alpha +p_{\overline{\theta }}^{}{}_{}{}^{\dot{\alpha }}\dot{\overline{\theta }}_{\dot{\alpha }}+p_{\mathrm{\Lambda }n}\dot{\mathrm{\Lambda }}^n+p_{\chi \alpha }\dot{\chi }^\alpha +p_{\overline{\chi }}^{}{}_{}{}^{\dot{\alpha }}\dot{\overline{\chi }}_{\dot{\alpha }}+\pi _\alpha \dot{\lambda }^\alpha +\pi ^{\dot{\alpha }}\dot{\lambda }_{\dot{\alpha }}+`$
$`\pi \dot{\lambda }+\pi _n\dot{\lambda }^n+\overline{𝒫}_\alpha \dot{C}^\alpha +\overline{𝒫}_{\dot{\alpha }}\dot{C}^{\dot{\alpha }}+\overline{𝒫}\dot{C}+\overline{𝒫}_n\dot{C}^n+\overline{C}_\alpha \dot{𝒫}^\alpha +\overline{C}_{\dot{\alpha }}\dot{𝒫}^{\dot{\alpha }}+\overline{C}\dot{𝒫}`$
$`+\overline{C}_n\dot{𝒫}^n\{\mathrm{\Psi }^{^{}},\mathrm{\Omega }^{^{}}\}_D),`$ (102)
where $`\mathrm{\Omega }^{^{}}`$ is given by Eq. (S0.Ex67) with the terms involving $`C^{1,2},\pi ^{1,2},\pi _{(1)}^{}{}_{}{}^{1,2}`$ omitted. In the course of the integration the standard change of variables (with unit Jacobian)
$$\pi ^{1,2}ϵ\pi ^{1,2},\overline{C}^{1,2}ϵ\overline{C}^{1,2},$$
(103)
followed by the limit $`ϵ0`$ has been used. Note that we do not see the compensating ghosts $`(C^{1,2},\overline{𝒫}^{1,2})`$ any more. The overcounting intrinsic to the sector $`(C^n,\overline{𝒫}_n)`$ is regulated now by the measure in Eq. (101).
In general, we can proceed on this way. However, this seems to break manifest Lorentz covariance. In particular we found that the following ansatz for $`\mathrm{\Psi }^{^{}}`$
$`\mathrm{\Psi }^{^{}}=\frac{1}{ϵ}\lambda ^+\overline{C}^{}+\frac{1}{ϵ}\lambda ^{}\overline{C}^++\frac{1}{ϵ}\mathrm{\Lambda }_i\overline{C}^i+\frac{1}{ϵ}\overline{C}_0\overline{\chi }^{\dot{1}}+\frac{1}{ϵ}\overline{C}_{\dot{0}}\chi ^1+\overline{𝒫}^i\lambda _i`$
$`+(\overline{𝒫}_0\frac{p^1ip^2}{\sqrt{2}p^{}}\overline{𝒫}_1)\lambda ^0+(\overline{𝒫}_{\dot{0}}\frac{p^1+ip^2}{\sqrt{2}p^{}}\overline{𝒫}_{\dot{1}})\lambda ^{\dot{0}}+\mathrm{\Psi }^{^{\prime \prime }},`$ (104)
where $`\mathrm{\Psi }^{^{\prime \prime }}`$ depends on $`(p^n,\theta ^\alpha ,p_{\theta }^{}{}_{\alpha }{}^{},\overline{\theta }^{\dot{\alpha }},p_{\overline{\theta }}^{}{}_{\dot{\alpha }}{}^{})`$ only and we switch to the light–cone notation $`\lambda ^n(\lambda ^\pm ,\lambda ^i)`$ (for fermions we write the indices explicitly $`\overline{C}_\alpha =(\overline{C}_0,\overline{C}_1)`$), reduces the integral to the standard path integral constructed with respect to the irreducible (noncovariant) subset of the Siegel constraints (1)
$$p_{\theta 0}\frac{p^1ip^2}{\sqrt{2}p^{}}p_{\theta 1}=0,p_{\overline{\theta }\dot{0}}\frac{p^1+ip^2}{\sqrt{2}p^{}}p_{\overline{\theta }\dot{1}}=0,p^2=0.$$
(105)
The explicit form of the latter is (we denote collectively $`z=(x^n,\theta ^1,\overline{\theta }^{\dot{1}})`$)
$$K(z_f,t_f;z_i,t_i)=\delta (\theta _{}^{1}{}_{f}{}^{}\theta _{}^{1}{}_{i}{}^{})\delta (\overline{\theta }_{}^{\dot{1}}{}_{f}{}^{}\overline{\theta }_{}^{\dot{1}}{}_{i}{}^{})K_0(x_f,t_f;x_i,t_i)$$
(106)
where $`K_0(x_f,t_f;x_i,t_i)`$ is the propagator of the massless spinless relativistic particle and $`(\theta _{}^{1}{}_{i}{}^{},\overline{\theta }_{}^{\dot{1}}{}_{i}{}^{},\theta _{}^{1}{}_{f}{}^{},\overline{\theta }_{}^{\dot{1}}{}_{f}{}^{})`$ denote the values of the fermions at the initial and final moments of time (boundary conditions).
Thus the path integral constructed above can be viewed as a formal covariantization of the propagator (106) characterizing the Siegel superparticle. Beautifully enough, this can be done with a finite number of ghost variables.
6. Discussion
In this article we have studied an alternative to the harmonic superspace approach, the latter seems to be the only method for quantizing infinitely reducible first class constraints currently available. The basic advantage of the novel technique is the existence of an explicit Lagrangian formulation and the validity of the standard spin–statistics relations for all the variables involved. In contrast to the harmonic superspace approach, where one first extracts linearly independent components from originally reducible constraints and then quantizes the resulting irreducible theory, the infinite reducibility of constraints is effectively canceled by that coming from the sector of auxiliary variables. Both methods, however, correlate well yielding a theory of rank two after $`BRST`$ quantization.
Turning to possible further developments, one expects the treatment of the $`ABCD`$ model along similar lines to be a natural next step. As has been mentioned in the Introduction, however, a proof of the equivalence of the $`ABCD`$ superparticle to a conventional model does not treat all constraints on equal footing. In view of this fact, the stringy extension seems to be preferable. Then, as was recently marked by Berkovits , a naive generalization of the present scheme to the superstring case faces the zero mode problem and, hence, deserves further investigation. We suspect, however, the latter point to be a technical difficulty rather than an ideological one. Another interesting point is to make use of the present approach to test an earlier quantization proposal by Kallosh (see also related works ,). The infinite proliferation of ghosts has been truncated there by imposing appropriate conditions on the ghosts variables, the latter involving specific (covariant) projectors. The phase space in our method is valid for the construction of such projectors (see also Ref. ) and the possibility to truncate the infinite ghost tower following Kallosh’s approach at the very second step seems to be tempting.
Acknowledgments
A.G. thanks Igor Bandos for useful discussions.
Appendix
In this Appendix we prove the equivalence of the last of Eqs. (S0.Ex16) and the pair (33), provided other constraints from (S0.Ex16) hold. Some details related to the analysis of the constraint system in the light–cone frame are also given.
Given the vector equation
$$\begin{array}{ccc}\hfill 2\varphi \mathrm{\Lambda }^n+\omega p^ni\phi \sigma ^n\overline{\chi }+i\chi \sigma ^n\overline{\phi }=0,& & \end{array}$$
$`(A.1)`$
the multiplication by $`\mathrm{\Lambda }_n`$ gives
$$\begin{array}{ccc}\hfill \omega =0.& & \end{array}$$
$`(A.2)`$
Hence, the second term in $`(A.1)`$ can be omitted. Passing to light cone coordinates, one has
$$\begin{array}{ccc}\hfill 2\varphi \mathrm{\Lambda }^+i\phi \sigma ^+\overline{\chi }+i\chi \sigma ^+\overline{\phi }=0,& & \end{array}$$
$`(A.3)`$
$$\begin{array}{ccc}\hfill 2\varphi \mathrm{\Lambda }^{}i\phi \sigma ^{}\overline{\chi }+i\chi \sigma ^{}\overline{\phi }=0,& & \end{array}$$
$`(A.4)`$
$$\begin{array}{ccc}\hfill 2\varphi \mathrm{\Lambda }^ii\phi \sigma ^i\overline{\chi }+i\chi \sigma ^i\overline{\phi }=0,& & \end{array}$$
$`(A.5)`$
where the customary notation $`\mathrm{\Lambda }^\pm =\pm \frac{1}{\sqrt{2}}(\mathrm{\Lambda }^0\pm \mathrm{\Lambda }^3)`$ is used.
It is worth mentioning now that, given a light–like vector $`\mathrm{\Lambda }^2=2\mathrm{\Lambda }^+\mathrm{\Lambda }^{}+\mathrm{\Lambda }^i\mathrm{\Lambda }^i=0`$, the equation $`(\phi \sigma ^n)_{\dot{\alpha }}\mathrm{\Lambda }_n=0`$ contains only half (one) linearly independent components. Actually, taking a conventional set of $`\sigma `$–matrices in $`R^{1|3}`$ (see Ref. )
$$\begin{array}{ccc}\{\sigma _n,\stackrel{~}{\sigma }_m\}=2\eta _{nm},\eta _{nm}=diag(,+,+,+),\hfill & & \\ \sigma ^+=\sqrt{2}\left(\begin{array}{cc}0& 0\\ 0& 1\end{array}\right),\sigma ^{}=\sqrt{2}\left(\begin{array}{cc}1& 0\\ 0& 0\end{array}\right),\hfill & & \\ \mathrm{\Lambda }_n\sigma ^n=\sqrt{2}\left(\begin{array}{cc}\mathrm{\Lambda }^+& \frac{\mathrm{\Lambda }^1i\mathrm{\Lambda }^2}{\sqrt{2}}\\ \frac{\mathrm{\Lambda }^1+i\mathrm{\Lambda }^2}{\sqrt{2}}& \mathrm{\Lambda }^{}\end{array}\right),\mathrm{\Lambda }_n\stackrel{~}{\sigma }^n=\sqrt{2}\left(\begin{array}{cc}\mathrm{\Lambda }^{}& \frac{\mathrm{\Lambda }^1i\mathrm{\Lambda }^2}{\sqrt{2}}\\ \frac{\mathrm{\Lambda }^1+i\mathrm{\Lambda }^2}{\sqrt{2}}& \mathrm{\Lambda }^+\end{array}\right),\hfill & & \end{array}$$
$`(A.6)`$
with $`\stackrel{~}{\sigma }^{n\dot{\alpha }\alpha }=ϵ^{\dot{\alpha }\dot{\beta }}ϵ^{\alpha \beta }\sigma _{n\beta \dot{\beta }}`$, one finds
$$\begin{array}{ccc}\hfill (\phi \sigma ^n)_{\dot{\alpha }}\mathrm{\Lambda }_n=0\{\begin{array}{cc}\phi ^0\mathrm{\Lambda }^++\phi ^1\frac{(\mathrm{\Lambda }^1+i\mathrm{\Lambda }^2)}{\sqrt{2}}=0\hfill & \\ \phi ^0\frac{(\mathrm{\Lambda }^1i\mathrm{\Lambda }^2)}{\sqrt{2}}+\phi ^1\mathrm{\Lambda }^{}=0.\hfill & \end{array}& & \end{array}$$
$`(A.7)`$
Multiplying the first equation in $`(A.7)`$ by $`\frac{(\mathrm{\Lambda }^1i\mathrm{\Lambda }^2)}{\sqrt{2}}`$ one recovers the second one, provided the standard light–cone condition
$$\begin{array}{ccc}\hfill \mathrm{\Lambda }^+0& & \end{array}$$
$`(A.8)`$
is assumed.
With the use of the explicit representation of the $`\sigma `$–matrices chosen, the constraint system $`(A.3)`$$`(A.5)`$ simplifies to
$$\begin{array}{ccc}\hfill 2\varphi \mathrm{\Lambda }^++i\sqrt{2}\phi ^1\overline{\chi }^{\dot{1}}i\sqrt{2}\chi ^1\overline{\phi }^{\dot{1}}=0,& & \end{array}$$
$`(A.9)`$
$$\begin{array}{ccc}\hfill 2\varphi \mathrm{\Lambda }^{}+i\sqrt{2}\phi ^0\overline{\chi }^{\dot{0}}i\sqrt{2}\chi ^0\overline{\phi }^{\dot{0}}=0,& & \end{array}$$
$`(A.10)`$
$$\begin{array}{ccc}\hfill \varphi (\mathrm{\Lambda }^1+i\mathrm{\Lambda }^2)i\phi ^0\overline{\chi }^{\dot{1}}+i\chi ^0\overline{\phi }^{\dot{1}}=0,& & \end{array}$$
$`(A.11)`$
$$\begin{array}{ccc}\hfill \varphi (\mathrm{\Lambda }^1i\mathrm{\Lambda }^2)i\phi ^1\overline{\chi }^{\dot{0}}+i\chi ^1\overline{\phi }^{\dot{0}}=0.& & \end{array}$$
$`(A.12)`$
On account of Eq. $`(A.7)`$ (the same holds for $`\chi `$ and complex conjugates), the last three equations follow from $`(A.9)`$. Thus, there appears to be only one linearly independent component entering the original vector equation. The latter can be put into a covariant (scalar) form. Actually, applying the same light–cone technique to the equation
$$\begin{array}{ccc}\hfill 2\varphi i\phi \sigma ^n\overline{\chi }p_n+i\chi \sigma ^n\overline{\phi }p_n=0,& & \end{array}$$
$`(A.12)`$
one recovers precisely Eq. $`(A.9)`$. |
warning/0002/quant-ph0002082.html | ar5iv | text | # Quantum searching with continuous variables
\[
## Abstract
A fast quantum search algorithm for continuous variables is presented. The result is the quantum continuous variable analog of Grover’s algorithm originally proposed for qubits. A continuous variable analog of the Hadamard (i.e., Fourier transform) operation is used in conjunction with inversion about the average of quantum states to allow the approximate identification of an unknown quantum state in a way that gives a square-root speed-up over search algorithms using classical continuous variables. Also, we show that this quantum search algorithm is robust for a generalised Fourier transformation on continuous variables.
\]
Quantum systems can register and process information in ways that classical systems cannot. As a result, it is possible for quantum computers to perform certain computational tasks faster than any classical computer . It is becoming increasingly clear that at the heart of quantum computation lies two basic quantum phenomena, one is quantum interference and the other quantum entanglement. The real upsurge of interest in this field came after Shor’s remarkable discovery of an algorithm for factoring large numbers . Subsequently, a fast quantum search algorithm has been discovered by Grover , which takes $`O(\sqrt{N})`$ steps instead of $`N`$ steps to search an unmarked item in a unsorted list of $`N`$ entries. Later the optimality of this search algorithm was proved . In particular, it was shown that this search algorithm can use almost any unitary transformation on qubit states . Further, it has even been argued that this algorithm would work even without entanglement (though at a cost in resources) .
These algorithms are usually implemented on quantum systems whose observables have discrete spectra, such as a collection of two-level atoms, ions, or spin-$`\frac{1}{/}2`$ particles (called qubits). However, there are other classes of quantum systems whose observables form continuous spectra. So, it is important to know how these algorithms can be generalised for continuous quantum variables? With the recent advances in our ability to manipulate continuous quantum information in teleportation , error-correction codes and its feasibility of implementation using linear devices , it is natural to ask whether one can provide some quantum algorithms that might be implemented on a continuous variable quantum computer. In fact, the usefulness of quantum computation over continuous variables has been recently emphasised . It has been shown that universal quantum computation over continuous variable is not only possible, but could be effected using simple non-linear operations with coupling provided solely by linear operations . These operations form a universal set of quantum gates for continuous variables allowing ‘quantum floating point’ arithmetic. While discrete quantum computation can be thought of as the coherent manipulation of qubits, continuous quantum computation can be thought of as the manipulation of ‘qunats’, where the qunat (pronounced as ‘Q nat’) is the unit of continuous quantum information.
In this letter, we propose a fast quantum search algorithm with continuous variables. Here a continuous variable can be anything, e.g., position, momentum, energy (unbounded) or amplitudes of the electromagnetic field. With the help of the Fourier transform (viewed as an active operation) on a continuous basis state (analogous to the Hadamard transform in the case of qubits) and a suitably generalised inversion operator, we construct a search operator which can be implemented on a continuous variable quantum computer. The inversion operator requires the projection operator for continuous basis states which we discuss. We show that the application of the compound operator takes $`O(\sqrt{N})`$ iterations to search an unmarked item in a list of $`N`$ entries. Further, we generalise our quantum searching with continuous variables to one based on a generalised Fourier transformation and which still gives a square root reduction in the number of steps. This shows that the quantum search algorithm with continuous variables is robust to the choice of arbitrary Fourier transformations. We also discuss the robustness of the search algorithm if one uses strongly peaked normalisable states instead of ideal infinite-energy position eigenstates.
Here, we discuss how to perform a quantum search algorithm using continuous variables. First we need to map a conventional discrete search problem into a continuous variable context. Suppose we have a function $`f(k):K\{0,1\}`$ defined on a domain $`K`$ with $`kK=\{1,2,\mathrm{},N\}`$. This function has a non-zero value equal to $`1`$ for some element $`k=k_f`$ and is $`0`$ for all other elements in the set $`K`$. Our task is to discover the value of $`k_f`$ given the ability to apply the function $`f`$ to inputs or superpositions of inputs, and given no further information about the function $`f(k)`$. In order to implement this in a quantum computer with continuous variables we require a collection of $`n`$ qunats. The state vector of each qunat belongs to a Hilbert space of infinite dimension. Since we have an infinite number of basis states, we cannot map each basis state within the Hilbert space onto each entry in the set $`K`$. (This would be a many-to-one mapping.) One could avoid this problem by choosing a subspace of the full Hilbert space with $`N`$ disjoint regions of the spectrum. However, we do not discuss this approach in detail here.
Let us consider a collection of $`n`$ continuous variables whose Hilbert space is spanned by a basis of states $`|x=|x_1,x_2,\mathrm{},x_n`$, satisfying the orthogonality condition $`x|x^{}=\delta (x_1x_1^{})\mathrm{}\delta (x_nx_n^{})=\delta (xx^{})`$. For example, one can consider a compact region of the state space divided into $`N`$ equal subvolumes, each with measure $`\mathrm{\Delta }x^n`$, one for each member of the set $`K`$. Let $`x_f`$ be the centre of the subvolume corresponding to $`k_f`$. In the context of this continuous variable embedding, executing the function $`f`$ corresponds to adjoining an extra state to the system, originally in the state $`|0`$, and applying an operator $`U_f:|x|0|x|1`$ if $`x`$ belongs to the region corresponding to $`k_f`$ and $`|x|0|x|0`$ otherwise. Clearly, if one samples the region at random by applying the operator to a series of random points, it will take $`O(N)`$ calls of the operator to find $`k_f`$.
If one exploits the power of quantum superposition and entanglement, however, fewer function calls are required. Our approach is discussed below. Let us pick an initial state in the position basis such as $`|x_i=|x_1,x_2,\mathrm{},x_n_i`$ for a quantum computer with continuous spectrum at random. The final (target) state is given by $`|x_f=|x_1,x_2,\mathrm{},x_n_f`$. We need a suitable unitary operator, which can take the initial state to the final state. Just as we have the Hadamard transformation in discrete computation, one of the basic operations with continuous variables is the Fourier transformation between position and momentum variables in phase space. By defining the Fourier transformation as an active operation on $`n`$ qunat states $`|x`$ we can write it as
$$|x=\frac{1}{\sqrt{\pi ^n}}𝑑ye^{2ixy}|y,$$
(1)
where $`xy=x_1y_1+\mathrm{}+x_ny_n`$, $`|y=|y_1,y_2,\mathrm{},y_n`$ and both $`x`$ and $`y`$ are in the position basis. This has been used by one of the present authors in developing an error correction code for continuous variables. This Fourier transformation can be straightforwardly applied in physical situations. For example, when $`|x`$ represents quadrature eigenstate of a set of modes of the electromagnetic field, $`|x`$ is simply an eigenstate of the conjugate quadrature.
Suppose, we apply the unitary operator $``$ to a basis state $`|x_i`$, then the relative amplitude of finding the system in the target state $`|x_f`$ is $`x_f||x_i=_{fi}=e^{2ix_ix_f}/\sqrt{\pi ^n}`$. Therefore, the relative probability of finding the system in the final qunat states will be given by $`|_{fi}|^2=1/\pi ^n`$. Hence, we have to repeat the experiment at least $`1/|_{fi}|^2=\pi ^n`$ times to successfully obtain the state $`|x_f`$. Here, we prove that search algorithm based on continuous variable can take $`\sqrt{\pi ^n}`$ steps to reach the final state starting from an initial state. (Here, we may identify the number of entries $`N`$ with $`\pi ^n`$.)
The next operator we need is the unitary operator, which can invert the sign of a basis state $`|x`$. We can define the selective inversion operator for a continuous basis $`|x`$ as
$$I_x=12P_{\mathrm{\Delta }x},$$
(2)
where $`P_{\mathrm{\Delta }x}`$ is the projection operator for continuous variables. Unlike the discrete case we cannot define the projection operator for the basis $`|x`$ as $`P_x=|xx|`$, because the operator $`P_x`$ is an ill defined and it will not satisfy $`P_x^2=P_x`$. The correct projection operator for continuous variables is defined as
$$P_{\mathrm{\Delta }x}=_{x_0\mathrm{\Delta }x/2}^{x_0+\mathrm{\Delta }x/2}𝑑x^{}|x^{}x^{}|.$$
(3)
The reason for this definition is that we cannot project an arbitrary state which is represented in terms of continuous basis state onto a point to get the exact eigenvalue. There will be always a spread within an interval. We can only project a state around $`x_0`$ to a selectivity $`\mathrm{\Delta }x`$ of the measuring apparatus. It is not possible to design a device to make a perfectly selective measurement of a continuous variable. The interval $`[x_1,x_2]`$ cannot be narrowed down, because it will always contains an infinite number of eigenvalues . Thus, if we have a wave packet the effect of projection is to truncate it around $`x_0`$ within an interval $`\mathrm{\Delta }x`$. This operator satisfies $`P_{\mathrm{\Delta }x}^2=P_{\mathrm{\Delta }x}`$ and $`P_{\mathrm{\Delta }x}|x=|x`$ as expected. With the help of the above inversion operator we can construct a compound search operator $`𝒞`$ defined as
$$𝒞=I_{x_i}^{}I_{x_f}.$$
(4)
It may be remarked that the selective inversion of the target state $`|x_f`$ can be achieved by attaching an ancilla qunat and considering the quantum XOR circuit for continuous variables . If a quantum circuit exists that transforms $`|x|a|x|f(x)+a`$, then by choosing the ancilla state $`|a=|\pi /2=𝑑ye^{i\pi y}|y/\sqrt{\pi ^n}`$ we can selectively invert the state $`|x`$ for which $`f(x)=1`$, i.e., $`|x|\pi /2|x|\pi /2`$.
Let us define a state $`|\stackrel{~}{x}_f^{}|x_f`$. We can show that the operator $`𝒞`$ can preserve the two-dimensional subspace spanned by the states $`|x_i`$ and $`|\stackrel{~}{x}_f`$. First, we show the action of $`𝒞`$ on $`|x_i`$. This can be expressed as
$`𝒞|x_i=|x_i4P_{\mathrm{\Delta }x_i}^{}P_{\mathrm{\Delta }x_f}|x_i+2^{}P_{\mathrm{\Delta }x_f}|x_i,`$ (5)
where $`P_{\mathrm{\Delta }x_i}=_{x_{i1}}^{x_{i2}}dx_i^{}|x^{}_i_ix^{}|`$, ($`x_{i1}=x_0\mathrm{\Delta }x_i/2`$, $`x_{i2}=x_0+\mathrm{\Delta }x_i/2`$ ) and likewise for $`P_{\mathrm{\Delta }x_f}`$. Using these facts we can simplify the above equation to
$`𝒞|x_i=(1{\displaystyle \frac{4}{\pi ^n}})|x_i+{\displaystyle \frac{2}{\sqrt{\pi ^n}}}{\displaystyle _{x_{f1}}^{x_{f2}}}𝑑x_f^{}e^{2ix_ix_f^{}}^{}|x_f^{}.`$ (6)
Similarly, we can evaluate the action of $`𝒞`$ on $`|\stackrel{~}{x}_f`$. It is given by
$`𝒞|\stackrel{~}{x}_f=|\stackrel{~}{x}_f{\displaystyle \frac{2}{\sqrt{\pi ^n}}}{\displaystyle _{x_{i1}}^{x_{i2}}}𝑑x_i^{}e^{2ix_i^{}x_f}|x_i^{}.`$ (7)
Thus, the operator $`𝒞`$ creates superpositions of two qunat states just as Grover’s operator creates superpositions of two qubit states. Once we understand the action of $`𝒞`$ on qunats we can obtain the total number of steps required in reaching the target state. Here, we use geometric structures from the projective Hilbert space of a quantum system to obtain the number of steps in the quantum searching. The projective Hilbert space admits a natural measure of distance called Fubini-Study distance . This measures the shortest distance between any two (not necessarily normalized) states $`|\psi _1`$ and $`|\psi _2`$ whose projections on $`𝒫`$ are $`\mathrm{\Pi }(\psi _1)`$ and $`\mathrm{\Pi }(\psi _2)`$, respectively. This can be defined as
$$d^2(|\psi _1,|\psi _2)=4\left(1\left|\frac{\psi _1}{||\psi _1||}|\frac{\psi _2}{||\psi _2||}\right|^2\right).$$
(8)
Here, the vectors $`|\psi _1`$ and $`|\psi _2`$ can be quantum states over continuous variables or over discrete variables. For unnormalisable states the above definition still works provided it is understood that the norm of the states can be made finite. In dealing with position eigenstates we can imagine that either the particle is moving in a finite space so that position eigenstates do not diverge or one can use normalisable states having strong peaks around some value of the position axis.
During the quantum searching with continuous variables we want to reach a state $`|\stackrel{~}{x}_f`$ from an initial state $`|x_i`$. This means we have to travel a shortest distance between these states which is given by $`d^2(|x_i,|\stackrel{~}{x}_f)=4(11/\pi ^n)`$. One application of the operator $`𝒞`$ creates a state $`|x_i^{(1)}=𝒞|x_i`$. We calculate the shortest distance between the resulting state $`|x_i^{(1)}`$ and the initial state $`|x_i`$. We note that the overlap of these states is given by $`x_i|𝒞|x_i=(14/\pi ^n)x_i|x_i+2\mathrm{\Delta }x_f/\pi ^n`$. For large database search $`N=\pi ^n`$ is very large and if we assume that the measuring device has a narrow selectivity, then $`\mathrm{\Delta }x_f`$ is also small. Hence, we can neglect the second term in the overlap (as it is a product of two small terms). Also note that term $`x_i|x_i`$ is not normalised but nevertheless it cancels out in during calculation. With this idea in mind we can evaluate the shortest distance between these states which is given by $`d^2(|x_i,|x_i^{(1)})32/\pi ^n`$. Thus in one application of the search operator $`𝒞`$ we can move the initial basis a shortest distance $`O(1/\sqrt{\pi ^n})`$. Therefore, to travel the full distance on the quantum state space we need $`N_s`$ number of steps given by
$$N_s=\frac{d(|x_i,|\stackrel{~}{x}_f)}{d(|x_i,|x_i^{(1)})}O(\sqrt{\pi ^n}).$$
(9)
This shows that a quantum computer based on qunats can take $`O(\sqrt{\pi ^n})`$ applications of $`𝒞`$ to reach the target state which otherwise would have taken $`O(\pi ^n)`$ number of steps by the application of $``$ on $`|x_i`$. Because the state is moving along a geodesic each application of $`𝒞`$ rotates the initial state in the right direction. This is our quantum search algorithm with continuous variables.
Instead of position eigenkets one can use strongly peaked normalisable state such as
$$|r_i=\frac{1}{(2\pi ϵ)^{n/4}}𝑑x\mathrm{exp}[\frac{(xx_i)^2}{4ϵ^2}]|x.$$
(10)
When $`ϵ0`$, the state $`|r_i`$ becomes a position eigenstate $`|x_i`$. Our algorithm can be practically implemented with such states. One can see that the action of the search operator $`𝒞`$ on $`|r_i`$ gives
$`𝒞|r_i`$ $`=`$ $`(1{\displaystyle \frac{4}{\pi ^n}})|r_i+{\displaystyle \frac{2}{\sqrt{\pi ^n}}}{\displaystyle \frac{1}{(2\pi ϵ)^{n/4}}}`$ (12)
$`\times {\displaystyle }dx{\displaystyle _{x_{f1}}^{x_{f2}}}dx_f^{}e^{(xx_i)^2/4ϵ^2+2ixx_f^{}}^{}|x_f^{}.`$
From the above formula one can see that a single application of search operator moves the initial state $`|r_i`$ a distance given by $`d^2(|r_i,|r_i^{(1)})48/\pi ^n`$. Now, if one defines the target state as
$$|r_f=\frac{1}{(2\pi ϵ)^{n/4}}𝑑x\mathrm{exp}[\frac{(xx_f)^2}{4ϵ^2}]|x,$$
(13)
then one can check that the total (shortest) distance between the initial state $`|r_i`$ and the desired state $`|\stackrel{~}{r}_f=^{}|r_f`$ is $`4[14O(ϵ^2)/\pi ^n]`$. Hence, by using (9) the total number of steps required to reach the target state is $`N_s=O(\sqrt{\pi ^n})`$.
Now we show that the quantum search algorithm over continuous variables is robust to some extent. Instead of the Fourier transform $``$ if we replace it by a generalised Fourier transform (GFT) in the search operator $`𝒞`$, still the algorithm works, i.e., we do get a square root reduction in the number of steps. We define a generalised Fourier transform as an active operation in the position basis $`|x`$ as
$`^{(\theta )}|x=({\displaystyle \frac{i}{\pi \mathrm{sin}\theta }})^{n/2}`$ (14)
$`\times {\displaystyle }dyexp[{\displaystyle \frac{i}{\mathrm{sin}\theta }}[(x^2+y^2)\mathrm{cos}\theta 2xy]]|y.`$ (15)
The GFT with a flexible angle $`\theta `$ gives a physical change of the basis $`|x`$ by any desired amount . Here, it should be mentioned that $`\theta >\mathrm{arcsin}(1/\pi )`$, since for smaller values of $`\theta `$ the assumption $`\mathrm{sin}^n\theta 1/\pi ^n`$ does not hold. The GFT for $`\theta =2\pi m`$, $`m`$ being an integer, corresponds to no change of basis. The GFT for $`\theta =\pi /2`$ corresponds to the Fourier transform defined in (1) (up to a constant phase shift equal to $`n\pi /4`$, $`n`$ being the number of qunats). If we apply GFT to an initial basis $`|x_i`$ then by probability rules of quantum theory we have to perform at least $`O[(\pi \mathrm{sin}\theta )^n]`$ number of trials to reach a target state $`|x_f`$. We will prove that the generalised search operator acting on continuous variables will take $`O[\sqrt{(\pi \mathrm{sin}\theta )^n}]`$ steps to reach the final state.
The search operator with this GFT takes the form
$$𝒞^{(\theta )}=I_{x_i}_{}^{(\theta )}{}_{}{}^{}I_{x_f}^{(\theta )}.$$
(16)
We can see that the action of the generalised search operator on the initial state $`|x_i`$ is given by
$`𝒞^{(\theta )}|x_i=\left(1{\displaystyle \frac{4}{(\pi \mathrm{sin}\theta )^n}}\right)|x_i+2\sqrt{{\displaystyle \frac{i^n}{(\pi \mathrm{sin}\theta )^n}}}`$ (17)
$`\times {\displaystyle _{x_{f1}}^{x_{f2}}}dx_f^{}\mathrm{exp}\{{\displaystyle \frac{i}{\mathrm{sin}\theta }}[(x_i^2+x_f^2)\mathrm{cos}\theta 2x_ix_f]\}^{}|x_f^{}.`$ (18)
Similarly, the action of the generalised search operator $`𝒞^{(\theta )}`$ on $`|\stackrel{~}{x}_f`$ can be calculated. It is given by
$`𝒞^{(\theta )}|\stackrel{~}{x}_f=|\stackrel{~}{x}_f2\sqrt{{\displaystyle \frac{(i)^n}{(\pi \mathrm{sin}\theta )^n}}}`$ (19)
$`\times {\displaystyle _{x_{i1}}^{x_{i2}}}dx_i^{}\mathrm{exp}\left\{{\displaystyle \frac{i}{\mathrm{sin}\theta }}[(x_i^2+x_f^2)\mathrm{cos}\theta 2x_i^{}.x_f]\right\}|x_i^{}.`$ (20)
It can be seen that the generalised search operator creates linear superposition of qunat states in the search process. Now, we can calculate the Fubini-Study distances to know how many steps are needed to reach the target state. The shortest distance between the states $`|x_i`$ and $`|\stackrel{~}{x}_f`$ is $`d^2(|x_i,|\stackrel{~}{x}_f)=4[11/(\pi \mathrm{sin}\theta )^n]`$. Notice that single application of the search operator $`𝒞^{(\theta )}`$ moves the initial state by a distance given by $`d^2(|x_i,|x_i^{(1)})=(32/\pi \mathrm{sin}\theta )^n`$. Therefore, to travel a shortest distance $`d(|x_i,|\stackrel{~}{x}_f)`$ we need $`N_sO[\sqrt{(\pi \mathrm{sin}\theta )^n}]`$ number of steps.
Thus, using a generalised Fourier transform we have proved that there is a square root reduction in the number of steps working with continuous variables. As expected for an angle $`\theta =\pi /2`$ we get back the original result with the search operator $`𝒞`$. This result is similar to the recent result of Grover , where the search algorithm for qubits has been generalised for arbitrary unitary transformations.
In conclusion, we have for the first time provided an efficient algorithm such as quantum searching to be implemented on a quantum computer with continuous variables. The key elements in this generalisation are the Fourier transformation and inversion operators which constitute the search operator for qunats in an infinite dimensional Hilbert space. We find that a square root speed up is possible with quantum computers based on qunats. Also, the continuous variable search is possible with almost any Fourier transformation. This may be practically implemented for any large data base search using linear and non-linear optical devices with the role for qunats being played by electromagnetic fields. It may well be that for large data base searches it is beneficial to use continuous quantum variables.
AKP and SLB acknowledge financial support from EPSRC. |
warning/0002/astro-ph0002204.html | ar5iv | text | # OBSERVED SMOOTH ENERGY IS ANTHROPICALLY EVEN MORE LIKELY AS QUINTESSENCE THAN AS COSMOLOGICAL CONSTANT
## 1 A Cosmological Constant or Quintessence?
Absent a symmetry principle protecting its value, no theoretical reason for making the cosmological constant zero or small has been found. Inflation makes the universe flat, so that, at present, the vacuum or smooth energy density $`\mathrm{\Omega }_{Q0}=1\mathrm{\Omega }_{m0}<1`$, is $`10^{120}`$ times smaller than would be expected on current particle theories. To explain this small but non-vanishing present value, a dynamic vacuum energy, quintessence, has been invoked, which obeys the equation of state $`w_QP/\rho <0`$. (The limiting case, $`w_Q=1`$, a static vacuum energy or Cosmological Constant, is homogeneous on all scales.)
Accepting this small but non-vanishing value for static or dynamic vacuum energy, the Cosmic Coincidence problem now becomes pressing: Why do we live when the clustered matter density $`\mathrm{\Omega }(a)`$, which is diluting as $`a^3`$ with cosmic scale $`a`$, is just now comparable to the static vacuum energy or present value of the smooth energy:
$$u_0^3\mathrm{\Omega }_{Q0}/\mathrm{\Omega }_{m0}2.$$
The observational evidence is for a flat, low-density universe:
(1) $`\mathrm{\Omega }_m+\mathrm{\Omega }_Q=1\pm 0.2`$
(Location of first Doppler peak in the CBR anisotroy at $`l200`$);
(2) $`\mathrm{\Omega }_{m0}=0.3\pm 0.05`$. (Slow evolution of rich clusters, mass power spectrum, CBR anisotropy, cosmic flows);
(3) $`\mathrm{\Omega }_{Q0}=1\mathrm{\Omega }_{m0}2/3`$ (curvature in SNIa Hubble diagram, dynamic age,height of first Doppler peak, cluster evolution).
Of these, the SNIa evidence is most subject to systematic errors due to precursor intrinsic evolution and the possibilty of grey dust extinction.
The combined data implies a flat, low-density universe with $`\mathrm{\Omega }_{m0}1/3`$, with negative pressure $`1w_Q1/2`$. In this paper, we use the evolution of large-scale structure to distinguish the two limiting cases:
LCDM: Cosmological constant: $`w_Q=1,n_Q3(1+w_Q)=0\mathrm{\Omega }_\mathrm{\Lambda }=2/3`$
QCDM: Quintessence: $`w_Q=1/2,n_Q=3/2,\mathrm{\Omega }_{Q0}=1/3`$ .
## 2 Evolution of a Low Density Flat Universe
The Friedmann equation in a flat universe with clustered matter and smooth energy density is
$$H^2(x)(\dot{a}/a)^2=(8\pi G/3)(\rho _m+\rho _Q),$$
or, in units of $`\rho _{cr}(x)=3H^2(x)/8\pi G`$, $`1=\mathrm{\Omega }_m(x)+\mathrm{\Omega }_Q(x),`$ where the reciprocal scale factor $`xa_0/a1+z\mathrm{}`$ in the far past, $`0`$ in the far future.
With the EOS $`wP/\rho `$, different kinds of energy density dilute at different rates $`\rho a^n,n3(1+w)`$, and contribute to the deceleration at different rates $`(1+3w)/2`$ shown in the table:
The expansion rate in present Hubble units is
$$E(x)H(x)/H_0=(\mathrm{\Omega }_{m0}x^3+(1\mathrm{\Omega }_{m0})x_Q^n)^{1/2}.$$
The Friedmann equation has an unstable fixed point in the far past and a stable attractor in the far future. (Note the tacit application of the anthropic principle: Why does our universe expand, rather than contract?)
The second Friedmann equation is $`\ddot{a}a/\dot{a}^2=(1+3w_Q\mathrm{\Omega }_Q)/2`$. The ratio of smooth energy to matter energy, $`\mathrm{\Omega }_Q/\mathrm{\Omega }_mu^3=u_0^3x^{3w_Q}`$, where $`\mathrm{\Omega }_{Q0}/\mathrm{\Omega }_{m0}u_0^32`$ is the present ratio. As shown by the inflection points in the middle curves of the figure, for fixed $`\mathrm{\Omega }_{Q0}/\mathrm{\Omega }_{m0}`$, QCDM (upper middle curve) expands faster than LCDM (lower middle curve), but begins accelerating only at the present epoch. The top and bottom curves refer respectively to a De Sitter universe ($`\mathrm{\Omega }_m=0`$), which is always accelerating, and an SCDM universe ($`\mathrm{\Omega }_m=1`$), which is always decelerating.
As summarized in the table below, quintessence dominance begins 3.6 Gyr earlier and more gradually than cosmological constant dominance. (In this table, the deceleration $`q(x)\ddot{a}/aH_0^2`$ is measured in present Hubble units.) The recent lookback time
$$H_0t_L(z)=z(1+q_0)z^2+\mathrm{},z<1,$$
where $`q_0=0`$ for QCDM and $`=1/2`$ for LCDM.
The density ratio $`u^3(a)\mathrm{\Omega }_Q/\mathrm{\Omega }_m=u_0^3x^{3w_Q}`$, increases as the matter density decreases. The matter-smooth energy transition $`\mathrm{\Omega }_Q/\mathrm{\Omega }_m=1`$ took place only recently at $`x^{w_Q}=u_0`$ or at $`x=1+z=u_0^2=1.5874`$ for QCDM and, even later, at $`x=1+z=u_0=1.260`$ for LCDM. Because, for the same value of $`u_0`$, a matter-QCDM freeze-out would take place earlier and more slowly than a matter-LCDM freeze-out, it imposes a stronger constraint on structure evolution. To permit evolution to the same present structure, QCDM would require a smaller value of $`\mathrm{\Omega }_{Q0}/\mathrm{\Omega }_{m0}`$ than does LCDM.
## 3 Growth of Large Scale Structure
The background density for large-scale structure formation is overwhelmingly Cold Dark Matter (CDM), consisting of clustered matter $`\mathrm{\Omega }_m`$ and smooth energy or quintessence $`\mathrm{\Omega }_Q`$. Baryons, contributing only a fraction to $`\mathrm{\Omega }_m`$, collapse after the CDM and, particularly in small systems, produce the large overdensities that we see.
Structure formation begins and ends with matter dominance, and is characterized by two scales: The horizon scale at the first cross-over, from radiation to matter dominance, determines the power spectrum $`P(k,a)`$, which is presently characterized by a scale factor $`\mathrm{\Gamma }=\mathrm{\Omega }_{m0}h=0.25\pm 0.05`$. The horizon scale at the second cross-over, from matter to smooth energy, determines a second scale factor, which for quintessence, is $`\mathrm{\Gamma }_Q`$ at $`130Mpc`$, the scale of voids, superclusters. A cosmological constant is smooth at all scales.
Quasars formed as far back as $`z5`$, galaxies at $`z6.7`$, ionizing sources at $`z=(1030)`$. The formation of any such structures, already sets an upper bound $`x<30`$ or $`(\mathrm{\Omega }_\mathrm{\Lambda }/\mathrm{\Omega }_{m0})<1000,\mathrm{\Omega }_{Q0}<30`$, for any structure to have formed. A much stronger upper bound, $`u_0<5`$, is set by when typical galaxies form i.e. by using the observed LSS, not to fix $`\mathrm{\Omega }_\mathrm{\Lambda }`$ or $`u_0^3`$, but to estimate the probability of our observing this ratio $`\mathrm{\Omega }_{Q0}/\mathrm{\Omega }_{m0}`$ at the present epoch.
For LCDM, Martel et al and Garriga et al calculate the asymptotic mass fraction that ultimately collapses into galaxies to be
$$f_{c,\mathrm{}}=\text{erfc}(\beta ^{1/2}),$$
remarkably a broad function of only $`\beta \delta _{i,c}^2/2(\sigma _i)^2`$, where $`\sigma _i^2=(1.72.3)/(1+z_i)`$ is the variance of the mass power spectrum and $`\delta _{i,c}`$ is the minimum density contrast which will make an ultimately bound perturbation. This minimum density contrast grows with scale factor $`a`$, and is approximately unity at recombination. Thus, except for a numerical factor of order unity , $`\delta _{i,c}x/(1+z_i)`$, the freeze-out projected back to recombination. Both numerator and denominator in $`\beta `$ refer to the time of recombination, but this initial time or red-shift cancels out in the quotient.
## 4 $`\mathrm{\Omega }_Q\mathrm{\Omega }_m`$ is Quite Likely for Our Universe
For a cosmological constant, an anthropic argument has already been given , assuming a universe of subuniverses with all possible values for the vacuum energy $`\rho _V`$ or $`\mathrm{\Omega }_\mathrm{\Lambda }`$. In each of these subuniverses, the probability for habitable galaxies to have emerged before the present epoch, is a function of $`\mathrm{\Omega }_\mathrm{\Lambda }`$ or the present ratio $`\mathrm{\Omega }_\mathrm{\Lambda }/\mathrm{\Omega }_{m0}`$
$$𝒫(\rho _V)(\text{prior distribution in}\rho _V)\times (\text{asymptotic mass fraction}f_{c,\mathrm{}}).$$
MSW, assuming nothing about initial conditions, assume a prior flat in $`\mathrm{\Omega }_\mathrm{\Lambda }`$. GLV argue that the prior should be determined by a theory of initial conditions and is not flat for most theories.
Following MSW, we assume a flat prior, so that the differential probability $`𝒫`$ for our being here to observe a value $`\rho _V`$ in our universe is simply proportional to the asymptotic collapsed mass fraction for this $`\rho _V`$. For LCDM,
$$\delta _{i,c}=1.1337u_0/(1+z_i),1.1337=(27/2)^{2/3}/5.$$
As function of the ratio $`\mathrm{\Omega }_\mathrm{\Lambda }/\mathrm{\Omega }_{m0}=u_0^3`$, the LCDM probability distribution has a broad peak about $`u_0^31230`$. The value observed in our universe $`u_0^32`$ has reasonable probability $`410\%`$.
This argument for LCDM ($`w_Q=1`$) is easily extended to QCDM ($`w_Q=1/2`$). The variance of the power spectrum, $`\sigma ^2`$, is insensitive to $`w_Q`$ for $`w_Q<1/3`$ . For $`w_Q=1/2`$, the numerical factor in $`\delta _{i,c}`$ is the same as for $`w_Q=1`$, but $`x=u_0^2`$ in place of $`u_0`$, so that $`\delta _{i,c}=1.1337u_0^2/(1+z_i)`$. Thus $`\beta _{QCDM}(u_0)=\beta _{LCDM}(\sqrt{u_0})`$, so that the QCDM probability distribution now peaks at $`u_0^33.55.5`$. With QCDM, the probability for observing $`u_0^32`$ is now increased to about 50%. |
warning/0002/astro-ph0002505.html | ar5iv | text | # Time Lags in Compact Objects: Constraints on the Emission Models
## 1. Models for the Formation of X/$`\gamma `$-ray Spectra
X-ray and gamma-ray spectra of accreting black holes and neutron stars are deconvolved into (at least) two components: a soft component interpreted as emission from an optically thick accretion disc, and a hard tail associated with a hot (10–100 keV) “corona”. Reviews of the spectral properties of Galactic black hole candidates (GBHs) can be found in Gilfanov et al. (1995), Tanaka & Lewin (1995), Grebenev et al. (1993, 1997), Grove et al. (1997), and Poutanen (1998). X/$`\gamma `$-ray properties of radio-quiet active galactic nuclei (AGN) are reviewed by Zdziarski et al. (1997), Johnson et al. (1997), and Zdziarski (1999). Recent results on the broad-band spectra of accreting neutron stars are presented by Barret et al. (2000).
An amusing fact is that super-massive black holes in AGN, GBHs and accreting neutron stars in their hard (low) states (see Tanaka & Lewin 1995; Gilfanov et al. 1995 for the definition of the spectral states) show very similar X/$`\gamma `$-ray spectra (see Zdziarski 1999; Barret et al. 2000). Furthermore, properties of their rapid temporal variability are also similar (van der Klis 1995b; Wijnands & van der Klis 1999; Psaltis, Belloni & van der Klis 1999; Ford et al. 1999; Edelson & Nandra 1999; Chiang et al. 2000). All this argues for a common mechanism of the X-ray production in all these sources.
There are good reasons to believe that the main radiative mechanism for the production of the hard X-rays is Comptonization of soft photons (e.g., Shapiro, Lightman, & Eardley 1976; Sunyaev & Trümper 1979; Sunyaev & Titarchuk 1980). However, it is not completely clear what determines the observed spectral slopes. The geometry of the X-ray emitting region and the source of soft photons is still a matter of debate (see Svensson 1996; Poutanen 1998; Beloborodov 1999b; Wardzinski & Zdziarski 2000).
An important clue to our understanding of the X-ray production came from the discovery of Fe lines (at $`6.4`$ keV) and the hardening of the spectra above 10 keV in AGN (Pounds et al. 1990; Mushotzky, Done, & Pounds 1993; Nandra & Pounds 1994), Cygnus X-1 (e.g., Done et al. 1992; Gierliński et al. 1997), and neutron stars (e.g., Yoshida et al. 1993). These features are associated with the reflection of hard X-rays from cold material (Basko, Sunyaev, & Titarchuk 1974; George & Fabian 1991; Magdziarz & Zdziarski 1995; Poutanen, Nagendra, & Svensson 1996). These observations gave support to the so called two-phase accretion disc-corona models. In such models, X-rays are emitted by a hot rarified corona above the cold accretion disc (Haardt & Maraschi 1993; Haardt, Maraschi, & Ghisellini 1994; Stern et al. 1995; Poutanen & Svensson 1996). Hard X-rays from the corona, being reprocessed in the cold disc, produce the reflection hump as well as most of the seed soft photons that are subsequently Comptonized to produce the hard X-rays. This is the feedback mechanism. The geometry of the corona determines the feedback factor which in its turn determines the spectral slope of the escaping radiation. The temperature of the emitting plasma (or to be more exact, the Kompaneets $`y`$-parameter) is determined by the energy balance between heating (by magnetic reconnection?) and cooling (by Comptonization of soft photons).
Further support for the feedback models was recently given by Zdziarski, Lubinski, & Smith (1999) (see also Zdziarski 1999; Gilfanov, Churazov, & Revnivtsev 2000) who found a correlation between the amount of reflection ($`R\mathrm{\Omega }/(2\pi )`$, where $`\mathrm{\Omega }`$ is a solid angle subtended by cold material as viewed from the X-ray source) and the intrinsic photon spectral index, $`\mathrm{\Gamma }`$, of the hard X-ray component. Such a correlation can easily be explained if there is overlap between the hot corona and the cold disc (Poutanen, Krolik, & Ryde 1997). The further the cold disc penetrates into the corona, the larger is the cooling, the smaller is the temperature of the corona, the softer is the spectrum, and, finally, the larger is the amplitude of the reflection. The model, however, appears to have trouble giving reflection amplitudes above $`R_{\mathrm{max}}0.5`$ (if the coronal optical depth $`\tau _\mathrm{T}1`$, see Zdziarski et al. 1997) due to partial smearing of the reflection component by the hot corona.
Alternatively, the observed $`R\mathrm{\Gamma }`$ correlation can be reproduced by variations of the bulk velocity of the X/$`\gamma `$-ray emitting plasma (Beloborodov 1999a,b). If the emitting regions are sufficiently compact to produce electron-positron pairs, the pressure of the radiation reflected and reprocessed in the disc accelerates pairs to mildly relativistic velocities away from the disc. On the other hand, for proton dominated plasmas, a small anisotropy in the energy dissipation mechanism can result in the ejection of particles away or towards the disc. Ejection away from the disc reduces $`R`$ below 1 and leads to hard spectra, while ejection towards to the disc can result in apparent $`R>1`$ as is observed in some objects.
The physical possibility of the corona formation was studied by Galeev, Rosner, & Vaiana (1979). They showed that the magnetic fields, being amplified in the cold disc due to turbulent motions and differential rotation, do not have time to annihilate inside the disc on the inflow time scale. Instead, the field loops are expelled from the disc by buoyancy (the Parker instability) and they annihilate in the tenuous corona. Beloborodov (1999a) showed that the mechanism studied by Galeev et al. is able to produce a corona of limited luminosity which is $`h/r`$ (the ratio of the disc height to its radius) times smaller than the disc luminosity. By contrast, in some sources most of the energy escapes in the form of hard X-rays. Beloborodov also argued that the magneto-rotational instability (Velikhov 1959; Chandrasekhar 1960; Balbus & Hawley 1991) increases the rate of the magnetic field generation (as compared with the Galeev et al. model) thus producing an active magnetic corona where a large fraction of the gravitational energy can finally be dissipated in magnetic flares. These qualitative arguments were recently supported by numerical three-dimensional magnetohydrodynamical simulations of Miller & Stone (2000) who showed that about 25 % of the total energy dissipation can occur in the rarified corona.
An alternative to the magnetic corona is the hot disc model (Shapiro et al. 1976; Ichimaru 1977; Narayan, Mahadevan, & Quataert 1998; Zdziarski 1998; Esin et al. 1998) which is also able to explain the observed X/$`\gamma `$-ray spectra. In order to distinguish between the models, it would be helpful to compare the predictions of different models with the temporal variability data (see van der Klis 1995a,b and Cui 1999a for recent reviews). Unfortunately, most of the papers on the spectral models do not consider the temporal variability. On the other hand, most of the models designed to explain the temporal variability data do not pay enough attention to the emission processes and the physics of the spectral formation.
In this review, we will discuss the variability data keeping in mind recent advances in modelling broad-band X/$`\gamma `$-ray spectra of accreting black holes and neutron stars. Most attention will be paid to the time lags that can shed light on the mechanism of the X-ray production. Then we discuss simple phenomenological models that are able to explain some of the observational facts. After that, we switch to the physical models. In particular, the properties of the Comptonizing regions will be discussed. We will point out the flaws in models that do not consider the energy balance in the “Compton cloud”, and then discuss models that satisfy the energy conservation law and confront them with the available data.
## 2. Observing Time Lags in Accreting Black Holes and Neutron Stars
The standard temporal characteristics that are usually computed are the power-density spectra (PDS) in different energy channels, auto/cross-correlation functions (ACF/CCF), the time/phase lags between the variability in different energy channels, the coherence function, etc. Time lags have been studied by two methods, by constructing the CCF and by cross-spectral analysis (for details, see Lewin et al. 1988; van der Klis 1989; Nowak et al. 1999a).
### 2.1. Lags in Black Hole Sources
The observations of Cygnus X-1 from sounding rockets and HEAO 1 (Priedhorsky et al. 1979; Nolan et al. 1981) showed that the CCF between different energy channels peaks very close to zero lag (delay $`<40`$ ms), but it is slightly asymmetric. Similar asymmetry was found in the EXOSAT data by Page (1985) who claimed a $`6\pm 1`$ ms shift of the peak of the CCF between the 5-14 keV and the 2-5 keV bands. Recent RXTE observations clearly show asymmetries of the CCFs, which, however, peak within $`1`$ ms from zero lag (see Fig. 1). This suggests that the relation between the variation in the two bands are not simply a time delay. Asymmetry is also observed in the RXTE data of GX 339-4, where the CCFs are offset by $`<5`$ ms from zero (using the 2-5 and 10-40 keV bands, Smith & Liang 1999). The CCFs of AGN also display similar properties (e.g., Papadakis & Lawrence 1995; Lee et al. 1999).
The CCFs cannot be fitted with simple exponentials at any time scale. A reasonably good description of the CCF is in terms of a stretched exponential $`CCF(t)=\beta \mathrm{exp}[(|tt_0|/\tau )^\nu ]`$ (see the inset of Fig. 1), where the normalisation $`\beta 1`$, the time where the CCFs peak $`t_0<10^3`$ s, $`\nu 2/3`$ in the range $`|t|<0.3`$ s, and the time constant $`\tau `$ is different for rising and decaying part of the CCF. Such a behaviour is probably the result of self-similarity of the light curve. It is interesting to note that the ACF of gamma-ray bursts also show a similar stretched exponential behaviour (see Stern & Svensson 1996; Beloborodov 1999c).
Since the CCF does not show which frequencies contribute most to the observed lags, van der Klis et al. (1987) suggested to use instead the cross-spectrum for such an analysis.<sup>1</sup><sup>1</sup>1The cross-spectrum $`C(f)S^{}(f)H(f)`$, where $`S(f)`$ and $`H(f)`$ are the Fourier transforms of the light curves in the soft and hard energy channels, respectively. The phase lag, $`\delta \varphi (f)\mathrm{arg}[C(f)]`$, and the time lag, $`\delta t(f)\delta \varphi (f)/(2\pi f)`$. The lags are positive when hard photons are lagging the soft ones. In the hard state of Cyg X-1 observed by the Ginga satellite, the time lags between the variability in the 1.2-4.7 and 4.7-9.3 keV energy bands reached 0.1 s and had a strong Fourier-frequency dependence $`\delta t(f)f^1`$, i.e., the phase lag $`\delta \varphi (f)`$ const (Miyamoto et al. 1988; Miyamoto & Kitamoto 1989). Similar lags were observed in other GBHs, GX 339-4 and GS 2023+338 ($``$V404 Cyg), in their hard state (Miyamoto et al. 1992). The analysis of the RXTE data for Cyg X-1 (Nowak et al. 1999a), GX 339-4 (Nowak et al. 1999b), 1E 1740.7-2942 and GRS 1758-258 (Smith et al. 1997), and GS 1354-644 (Revnivtsev et al. 2000) confirming the general features seen in the Ginga data, showed more complicated behaviour of the phase lag spectra which have a number of shelves and breaks (see Fig. 2).
Grove et al. (1998) extended this analysis to higher energies with the data from CGRO/OSSE. The time lags at low frequencies reached 0.3 s between the 50-70 and the 70-100 keV photons in the light curve of GRO J0422+32 ($``$Nova Persei 1992). The breaks detected in the time lag spectrum at 0.1 Hz may be related to the quasi-periodic oscillation (QPO) observed at 0.23 Hz. The CGRO/BATSE data of Cyg X-1 (Crary et al. 1998), GRO J0422+32 and GRO J1719-24 ($``$Nova Oph 1993) (van der Hooft et al. 1999a,b) show very similar time lag spectra.
The time lags of GBHs in their soft state (when $`\mathrm{\Gamma }2.5`$) turn out to be somewhat different. Soft lags were observed between the 1.2-2.3 keV and the 2.3-4.6 keV bands in GX 339-4 (Miyamoto et al. 1991) and GS 1124-68 ($``$Nova Muscae 1991; see Miyamoto et al. 1993; Takizawa et al. 1997), while the higher energy photons were lagging the variability in the 2.3-4.6 keV band. In one observation of GS 1124-68, the variability in the 4.6-9.2 keV was the most advanced. The phase lag reached $`1`$ rad which is much larger than the lags in the broad band noise observed in the GBHs in their hard state. Rapid time variations were mostly due to the harder power-law component which is clearly seen from the rms amplitude. It is interesting that the lags seem to saturate above 10 keV. It is worth pointing out that the largest lags here are observed at the QPO frequency. On the other hand, the time lag spectrum of Cyg X-1 in the soft state looks quite similar to that in the hard state (Cui et al. 1997).
Time lags have been observed in GRS 1915+105 in the 67 mHz QPO by Cui (1999b) and in the broad-band noise and QPOs by Reig et al. (2000). The lags show a very complicated structure, sometimes changing signs from one harmonic to another, and the sign also depends on the frequency of the QPO. Wijnands, Homan, & van der Klis (1999) and Cui, Zhang, & Chen (2000) observed lags in the broad-band noise and the QPOs of XTE J1550-564.
### 2.2. Lags in Neutron Star Sources
Hard time lags in neutron star sources were discovered by Hasinger (1987) in the CCF of Cyg X-2 (comparing the 1-5 and the 5-17 keV bands) in its horizontal branch (HB, see Lewin et al. 1988 for definitions of branches). The CCF also had sinusoidal oscillations due to the QPO. The time lags showed anti-correlation with the QPO frequency, dropping from 4 ms to 1.5 ms when the QPO changed from 20 Hz to 50 Hz. Hasinger interpreted the lags as delays due to scattering (Comptonization) in the hot cloud and the anti-correlation as an indication of a change in the system size. Associating the QPO frequency with the Keplerian frequency at some radius gives the relation, $`\delta t(f)f_{QPO}^{2/3}`$, while the actual data are much better described by $`\delta tf_{QPO}^1`$, i.e. $`\delta \varphi =2\pi f\delta t=\mathrm{const}`$.
Using the cross-spectrum techniques van der Klis et al. (1987) confirmed the existence of $`3`$ ms hard lags in the 20-40 Hz QPO of Cyg X-2 (and GX 5-1) and discovered 8 ms soft lags in the low-frequency noise, which were interpreted as a softening of the spectrum during the shots that cause QPOs. These results were confirmed by Vaughan et al. (1994) who also showed (from the analysis of the Ginga data of GX 5-1 on the HB) that the time lags increase with photon energy.
In both Cyg X-2 and GX 5-1 in their normal branch, the lags at the $`5`$ Hz QPO showed energy dependence (Mitsuda & Dotani 1989; Vaughan et al. 1999) reaching $`\delta \varphi \pi `$ rad (i.e., $`\delta t0.2\mathrm{s}`$) for 10 keV photons vs 2 keV. At the same time, the rms amplitude of the QPO in Cyg X-2 had a minimum at 5 keV and in GX 5-1 it increased above 2.5 keV. This behaviour can be interpreted as a pivoting of the spectrum around 3-5 keV.
With the larger effective area of RXTE, Ford et al. (1999) and Olive & Barret (2000) discovered phase lags in the broad-band noise of three atoll sources, 4U0614+09, 4U1705-44, and 4U1728-34. These lags are very similar to those in GBHs like Cyg X-1 and GX 339-4, which tells us that the mechanism responsible for the lags does not depend on the presence or absence of the hard surface of the neutron star, magnetosphere, boundary layer, etc., but instead is a property of the accretion flow.
A number of neutron stars show kHz QPOs in their light curves as revealed by RXTE. Kaaret et al. (1999) find 25 $`\mu `$s soft lags between the 4-6 keV and the $`>`$ 9 keV photons in the 800 Hz QPO in the atoll source 4U1636-536. Analysing the 550 Hz oscillations of Aquila X-1, Ford et al. (1999) found soft lags, $`\delta \varphi 1`$ rad, between the 3-6 keV and the $`>`$ 6 keV photons. Similar lags were found in the accreting millisecond pulsar SAX J1808.4-3658 (Cui, Morgan, & Titarchuk 1998; Ford 2000). The lags in the 830 Hz QPO in 4U 1608-52 (Vaughan et al. 1997, 1998) reach 60 $`\mu `$s between the 5 and the 25 keV photons.
## 3. Phenomenological Energy Dependent Shot Noise Models
It is clear that the observed zoo of time lags cannot be explained by any single model. Different mechanisms should be involved in producing the lags at different Fourier frequencies, the lags in the QPOs and the coherent pulsations, and the lags in the broad-band noise.
Let us first consider the simplest possible model that produces time lags: a shot noise model (Terrell 1972), where shots (=flares) are uncorrelated with each other. We assume that the shot time profiles at different energies have the same shape, but slightly different time constants. As an example we take a shot profile at soft energies $`s(t)=[t/\tau ]^p\mathrm{exp}[t/\tau ]`$ and at hard energies $`h(t)=[t/(\eta \tau )]^p\mathrm{exp}[t/(\eta \tau )]`$, where $`t>0`$ is measured from the beginning of the shot and $`p`$ is positive. The Fourier transforms, $`S(f)`$ and $`H(f)`$, are $`S(f)\tau \mathrm{\Gamma }(p+1)/(1i2\pi f\tau )^{p+1}`$ and $`H(f)\eta \tau \mathrm{\Gamma }(p+1)/(1i2\pi f\eta \tau )^{p+1}`$. The PDSs are $`|S(f)|^2`$ and $`|H(f)|^2`$. For small frequencies, $`f1/(2\pi \tau )`$, the PDSs have a flat dependence on frequency, $`f^0`$, while for large frequencies, $`f1/(2\pi \tau )`$, the PDSs decay as $`f^{2(p+1)}`$. The power per logarithm of frequency (i.e., $`f\times PDS(f)`$) peaks for soft photons at $`f_{s,\mathrm{max}}=1/[2\pi \tau \sqrt{2p+1}]`$ and at $`f_{h,\mathrm{max}}=1/[\eta 2\pi \tau \sqrt{2p+1}]`$ for hard photons. The phase lags $`\delta \varphi (f)=(p+1)[\mathrm{arctan}(\eta 2\pi f\tau )\mathrm{arctan}(2\pi f\tau )]`$. The lag is positive when hard photons are lagging soft ones (i.e., for $`\eta >1`$). For small frequencies, $`\delta \varphi (f)`$ rises as $`(p+1)(\eta 1)2\pi f\tau `$, while for large frequencies, it decays as $`(p+1)(\eta 1)/(\eta 2\pi f\tau )`$. The lag reaches a maximum of $`\delta \varphi _{\mathrm{max}}=2(p+1)(\mathrm{arctan}\sqrt{\eta }\pi /4)`$ at $`f=1/(2\pi \tau \sqrt{\eta })`$ close to $`f_{s,\mathrm{max}}`$ and $`f_{h,\mathrm{max}}`$, the frequencies where $`f\times PDS_{s,h}(f)`$ peak.
One can also consider a modified shot noise model, where the shot time scales are distributed according to a power law, $`\rho (\tau )\tau ^p`$ between $`\tau _{\mathrm{min}}`$ and $`\tau _{\mathrm{max}}`$ (see, e.g., Miyamoto & Kitamoto 1989; Lochner, Swank, & Szymkowiak 1991), with the same ratio $`\eta `$. Physically this could correspond to, for example, the situation when flares of different durations appear at different radii from the central object (Poutanen & Fabian 1999a). A power-law distribution of $`\tau `$ assures that the PDS is also a power-law $`f^{(3p)}`$ (Lochner et al. 1991). If the flares are self-similar, then the phase lag will be constant $`\delta \varphi _{\mathrm{max}}`$ for $`ff_{\mathrm{min}}1/(2\pi \tau _{\mathrm{max}})`$ and $`ff_{\mathrm{max}}1/(2\pi \tau _{\mathrm{min}})`$, decay as $`1/f`$ at $`f>f_{\mathrm{max}}`$ and rise linearly at $`f<f_{\mathrm{min}}`$. The corresponding time lags are constant, $`\delta t_{\mathrm{max}}=2\pi \tau _{\mathrm{max}}\delta \varphi _{\mathrm{max}}`$, for $`f<f_{\mathrm{min}}`$, and decay approximately as $`1/f`$ between $`f_{\mathrm{min}}`$ and $`f_{\mathrm{max}}`$. Note that in this model the coherence function (Vaughan & Nowak 1997; Nowak et al. 1999a) is close to unity, since the light curves at different energies are almost perfectly synchronised.
If we assume that $`\eta >1`$, there are hard lags and the predicted behaviour of the time lags and coherence function is in a very good agreement with the observations of GBHs (Poutanen & Fabian 1999a). However, this model contradicts the CCF of Cyg X-1 (see Fig. 1 and Maccarone et al. 2000). The CCF becomes narrower at larger energies which requires the shots to be narrower at larger energies. If one, however, reverses the time profiles of the shots, so that they rise slower and decay faster (e.g., $`s(t)=(t/\tau )^p\mathrm{exp}(t/\tau ),t<0`$), and one assumes $`\eta <1`$ (i.e., hard shots are narrower), the CCFs and the time lags can be reproduced, simultaneously.
Much more complicated models which also account for lags in the QPOs sources were developed by Shibazaki et al. (1988). We just note here that if a signal consists of shots appearing almost periodically and if shots at different energies are shifted in time one against another (or, e.g., the minima are reached at the same time and the peaks are not), the phase lag has a very complicated dependency on frequency (e.g., changes sign from one harmonic to another) depending on the shot profiles.
## 4. Physical Mechanisms for Producing Lags
### 4.1. Static Compton Cloud Models
Since Comptonization is the most probable mechanism for X-ray production in compact objects, it is natural to attribute the time delays between hard and soft photons to this process. Hard photons are the result of more scattering and so emerge after, or lag behind, softer ones. Consider a static “Compton cloud” with fixed Thomson optical depth, $`\tau _\mathrm{T}`$, and electron temperature $`\mathrm{\Theta }=kT_e/m_ec^2`$. A soft seed photon of energy $`E_0`$ injected into the cloud increases its energy by a factor of $`A_1=1+4\mathrm{\Theta }+16\mathrm{\Theta }^2`$ on average after each scattering, so that after $`N`$-scatterings its energy $`E_N=A_1^NE_0`$. The photon mean free path is $`\lambda R/\mathrm{max}(1,\tau _\mathrm{T})`$ (where $`R`$ is the size of the X-ray producing region, and where we accounted for the fact that we are interested only in those photons that actually have undergone scatterings in the cloud). The time between successive scatterings is then $`t_c=R/(c\mathrm{max}[1,\tau _\mathrm{T}])`$, so the time needed to reach the energy $`E_N`$ is (Sunyaev & Titarchuk 1980; Payne 1980)
$$t_N=Nt_c=\frac{R/c}{\mathrm{max}(1,\tau _\mathrm{T})}\frac{\mathrm{ln}(E_N/E_0)}{\mathrm{ln}A_1},$$
which translates to $`t_N10^4\mathrm{s}`$ for $`kT_e50`$ keV, $`\tau _\mathrm{T}1`$, $`R=10`$ km, and $`E_N/E_010`$.
This model was criticised by Miyamoto et al. (1988), Miyamoto et al. (1991) and Vaughan et al. (1994). First, the large size of the cloud ($`10^310^5R_g`$, where $`R_g`$ is the Schwarzschild radius, $`2GM/c^2`$) is needed to produce large delays observed in GBHs and neutron stars. Such cloud is physically unrealistic, since most of the gravitational energy is dissipated within $`10R_g`$. Second, the lags predicted by the model are independent of the Fourier frequency (Miyamoto et al. 1988) contrary to the observed $`1/f`$ dependence. Third, it is assumed that the soft photons produce the variability, while the hot cloud is not variable. Observationally, it is well established that when a soft black body spectrum is observed in GBHs it is much less variable than the hard X-rays (e.g. Miyamoto et al. 1991), so that the hard X-ray variability is most probably intrinsic to the hot cloud itself.
Finally, we would like to point out that due to the requirement of energy conservation the whole concept of a static Compton cloud with a constant temperature is physically unrealistic. The total emitted X-ray luminosity (produced by Comptonization of soft radiation) is $`L_{\mathrm{tot}}=L_h+L_s`$, where $`L_h`$ is the heating rate in the hot cloud and $`L_s`$ is the luminosity of seed soft photons. For hard spectra (i.e. $`L_hL_s`$), $`L_{\mathrm{tot}}L_h`$. The total X-ray luminosity is thus a function of the heating rate only and it does not depend on the amount of seed soft photons. By changing $`L_s`$, one effectively changes the spectral slope of the emergent X-ray radiation which is a function of the Compton amplification factor $`AL_h/L_s`$. This results in the pivoting of the spectrum (see Poutanen 1998; Beloborodov 1999b) without a noticeable increase in $`L_{\mathrm{tot}}`$. The larger the $`L_s`$, the smaller the equilibrium temperature of the emitting electrons, and the softer the spectrum. By contrast, in static Compton cloud models no changes in the electron temperature are considered in reaction to the changes in the number of soft photons, violating thus the energy conservation law. In order to increase the emitted luminosity, one has to change the energy dissipation rate in the cloud, but then exactly these changes will be driving the variability.
Recently, Kazanas, Hua, & Titarchuk (1997) (see also Böttcher & Liang 1998; Hua, Kazanas, & Cui 1999) modified the Comptonization model. Instead of a homogeneous Compton cloud, the density profile, $`n(r)1/r`$, was assumed (in this case, one has equal optical depth per logarithm of radius). The variability is still driven by changing the rate of soft photon injection in the center of the cloud. Then larger radii produce lower frequency variability (filtering out high frequency signal) and larger lags, while the smaller radii produce higher frequency variability and smaller lags (see also Nowak et al. 1999a,c). This model solves only one of the problem mentioned above ($`1/f`$ time lag dependence), while the other problems remain unsolved. Another modification of the model was considered by Böttcher & Liang (1999) based on an earlier suggestion by Miyamoto et al. (1988). Here, small cold clouds are assumed to free-fall into the hot central cloud thus changing the input of soft photons. Again, this model has problems with energy conservation.
### 4.2. The Dynamic Compton Cloud
Miyamoto et al. (1988) pointed out that some modulation mechanism must be invoked to produce the strongly frequency-dependent time lags. Poutanen & Fabian (1999a,b) proposed a model where the time lags are produced by the evolution of the flare spectrum. They assumed that the energy dissipation varies in time. For a small (of the order of $`R_g`$) emitting region (ER) one can consider the spectral evolution as a sequence of steady-states as long as the characteristic time scale of variability is $`\tau R/c`$. Any changes in the amplification factor $`A`$ would cause spectral variability and, specifically, a continuous increase of $`A`$ with time during the course of the flare would cause a soft-to-hard spectral evolution producing hard time lags of the order of the flare time scale (see § 3.). Poutanen & Fabian (1999b) considered three mechanisms that can increase $`A`$.
(1) A flare starts in the background of soft photons. The X-ray spectrum is soft as long as $`L_h<L_{s,\mathrm{bkg}}`$. With increasing $`L_h`$, the soft photon input gets dominated by reprocessed photons. The spectral slope is then determined by the feedback parameter $`D1/A`$ (the fraction of $`L_h`$ returned to the ER after the reprocessing in the disc into soft seed photons, see Stern et al. 1995; Svensson 1996; Beloborodov 1999b) and the spectrum becomes hard.
(2) The dissipation is accompanied by pumping net momentum into the hot plasma of the emission region (Beloborodov 1999a,b). The resulting bulk velocity increases with increasing luminosity. It leads to a lower feedback and higher $`A`$ (if the velocity is directed away from the disc).
(3) The differential rotation of the footpoints of a magnetic loop at the disc surface causes a twisting and elevation of the loop (Romanova et al. 1998). The time scale of the evolution is of the order of the Keplerian time-scale. When the emission region moves away from the disc, the feedback decreases and $`A`$ increases (see Fig. 3).
In all these cases the spectral evolution proceeds from soft to hard. The hard time lags between energies $`E`$ and $`E_0`$ are $`\tau \mathrm{ln}(E/E_0)`$. If there is a distribution of time scales $`\tau `$ between, say, 1 ms and 0.3 s (e.g., Keplerian time scales at the radii between the innermost radius of the accretion disc and $`50R_g`$), the time lags $`1/f`$ at the characteristic frequencies of the variability (see § 3. and Fig. 2).
If $`\tau `$ a few light crossing time of the ER, the spectrum in the beginning of the flare is hard because of photon starvation (one needs a few $`R/c`$ to get reprocessed soft photons into the ER) and softens towards the end of the flare (Poutanen & Fabian 1999a; Malzac & Jourdain 2000) producing soft time lags. Observing the change of sign of the time lags at some Fourier frequency, $`f_{\mathrm{sgn}}`$, would determine the size of the ER that produces variability at these frequencies (e.g., for Cyg X-1, $`R<10R_g(30`$ Hz$`/f_{\mathrm{sgn}})`$).
### 4.3. Small Scale Spectral Transitions
The models considered above can explain the time lags in the broad-band noise. In some QPO sources the time lags show a very complicated behaviour (see § 2.) and may require different explanation. In radiation-hydrodynamic model (see, e.g., Lamb 1989; Miller & Lamb 1992), QPOs appear as a result of oscillations in the optical depth of the radial flow due to the radiation feedback from the neutron star surface. The resulting spectral pivoting produces phase lags and the increase of the rms amplitude variability above the pivoting point. This model, however, is not applicable to the QPOs and the lags in black holes sources (see, e.g., Takizawa et al. 1997) because of the absence of a hard surface.
The galactic microquasar GRS 1915+105 shows large amplitude oscillations with periods varying from $`<`$1 up to 100 s. These time scales are a few orders of magnitudes larger than Keplerian time scales and up to $`10^6`$ time larger than the light crossing time of the ER. The best candidate for producing the spectral variability that causes the lags in GRS 1915+105 is the oscillation of the inner radius of the cold disc on viscous time scales. Such oscillations are similar to the spectral transitions, but have smaller amplitude and occur at shorter time scales. Changes of the relative geometry of the hot corona and the cold disc (with or without changes of the total luminosity) cause spectral pivoting at a few keV (see, e.g., Poutanen, Krolik, & Ryde 1997; Esin et al. 1998). The fluxes below and above the pivot point oscillate then with a phase shift of $`\pi `$. The rms amplitude of the QPO increases with the energy. The phase lags between the energies above the pivot point can then be produced if the oscillations are time asymmetric (see Fig. 11 in Morgan, Remillard, & Greiner 1997; Vilhu & Nevalainen 1998).
Similar (but aperiodic) changes in the inner disc radius can be responsible for the broad-band variability observed at $`f<1`$ Hz in, e.g., Cyg X-1. Associated spectral changes can manifest themselves in time lags observed at these frequencies.
### 4.4. Delays due to Compton Reflection and Reprocessing
The spectra of accreting GBHs and neutron stars show signatures of Compton reflection (see § 1.). Some fraction of the X-ray photons can be reflected from the outer edge of a flared accretion disc, a wind from the companion, etc. Such a reflector acts as a low pass filter smearing out the high frequency variations and produces lags corresponding to the light travel time to the reflector only at lower frequencies. Such processes can explain the break in the time lag spectra observed in Cyg X-1 (see Fig. 2) and in other GBHs at $`f<1`$ Hz.
Reprocessed soft radiation which accompanies Compton reflection is emitted in the optical and UV spectral bands if the reprocessing occurs far away from the central X-ray source. The time delays can then be measured between the optical/UV and the X-ray radiation (e.g., Hynes et al. 1998). On the other hand, reprocessing in the vicinity of the X-ray emitting region, produces time delays of the order of the time scale of the spectral evolution in the hard X-ray band. This may be one reason for the observed soft lags in the soft state of GX 339-4 and GS 1124-68 (§ 2.1.).
### 4.5. Hot Spots on the Neutron Star Surface
Some neutron star sources show lags in their periodic oscillations at kHz frequencies. Ford et al. (1999) and Ford (2000) interpreted the soft lags in Aquila X-1 and in the accreting millisecond pulsar SAX 1808.4-3658 using a model of a rotating hot spot with a black body spectrum at the surface of a neutron star where the lags appear due to Doppler effects. The weak energy dependence of the rms amplitude of the oscillations reported by Cui et al. (1998) rules out the black body model for the spectrum used by Ford. Detailed analysis of the pulsations in the time domain by folding techniques by Revnivtsev (1999) revealed that the pulse profile is distorted at different energies, while the minima are reached at the same time (i.e., there are no lags in the normal meaning of this word).
## 5. Summary
Time lags and other temporal variability data provide strong constraints on the models of the X-ray production. It was demonstrated that static Compton cloud models are based on physically unrealistic assumptions. The models invoking spectral evolution of the flare spectrum can fit both the CCF and the time lag Fourier spectra only if (1) the energy dissipation rate increases slowly and decreases rapidly and (2) the flare spectrum evolves from soft to hard. If soft seed photons are produced by reprocessing the hard ones, the change of sign in the time lag spectrum is expected at high frequencies corresponding to the light crossing time of the emission region. The absence of such a change would put constraints on the size of the emitting region.
We also argued that the reflection of hard X-rays from the outer part of the accretion disc produces time delays that we already might have observed in GBHs. If so, the disc should be flared and the break in the time lag Fourier spectra then corresponds to the size of the accretion disc. Of course, such an interpretation is not unique. Alternatively, small scale spectral transitions (e.g., oscillations of the inner radius of the accretion disc at viscous time scales) might produce time lags observed at lower frequencies.
In the case of (quasi-) periodic oscillations from the neutron star sources, we argued that in order to reproduce both the time lags and the energy dependent rms amplitude, the spectrum of the hot spots should not be close to a black-body.
## ACKNOWLEDGEMENTS
This work was supported by the Swedish Natural Science Research Council and the Anna-Greta and Holger Crafoord Fund. I thank Katja Pottschmidt for providing the time lag Fourier spectra and the light curves of Cyg X-1 used in the calculations of the cross-correlation functions. I am grateful to Andrei Beloborodov and Roland Svensson for valuable comments.
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warning/0002/math0002049.html | ar5iv | text | # 𝐿²–Riemann–Roch inequalities for covering manifolds
## §1. Estimates of the spectrum distribution function
Let $`\stackrel{~}{M}`$ be a complex analytic manifold of complex dimension $`n`$ on which a discrete group $`\mathrm{\Gamma }`$ acts freely and properly discontinuously. Let $`X=\stackrel{~}{M}/\mathrm{\Gamma }`$ let $`\pi :\stackrel{~}{M}X`$ be the canonical projection. We assume $`\stackrel{~}{M}`$ paracompact so that $`\mathrm{\Gamma }`$ will be countable. Suppose we are given a holomorphic vector bundle $`F`$ on $`X`$ and take its pull-back $`\stackrel{~}{F}=\pi ^{}F`$, which is a $`\mathrm{\Gamma }`$ invariant bundle on $`\stackrel{~}{M}`$. We also fix a $`\mathrm{\Gamma }`$ invariant hermitian metric on $`\stackrel{~}{M}`$ and on $`\stackrel{~}{F}`$. We consider a relatively compact open set $`\mathrm{\Omega }X`$ and its preimage $`\stackrel{~}{\mathrm{\Omega }}=\pi ^1\mathrm{\Omega }`$; $`\mathrm{\Gamma }`$ acts on $`\stackrel{~}{\mathrm{\Omega }}`$ and $`\stackrel{~}{\mathrm{\Omega }}/\mathrm{\Gamma }=\mathrm{\Omega }`$. In general we will decorate by tildes the preimages of objects living on the quotient. Let $`U`$ be a fundamental domain of the action of $`\mathrm{\Gamma }`$ on $`\stackrel{~}{\mathrm{\Omega }}`$. This means that (see e.g. \[At\]): a) $`\stackrel{~}{\mathrm{\Omega }}`$ is covered by the translations of $`\overline{U}`$, b) different translations of $`U`$ have empty intersection and c) $`\overline{U}U`$ has zero measure (since $`\mathrm{\Omega }`$ is smooth). Since $`\mathrm{\Omega }`$ is relatively compact $`U`$ has the same property. Let us define the space of square integrable sections $`L^2(\stackrel{~}{\mathrm{\Omega }},\stackrel{~}{F})`$ with respect to a $`\mathrm{\Gamma }`$ invariant metric on $`\stackrel{~}{M}`$ (and its volume form) and a $`\mathrm{\Gamma }`$ invariant metric on $`\stackrel{~}{F}`$. Then $`L^2(U,\stackrel{~}{F})`$ is constructed with respect to the same. There is a unitary action of $`\mathrm{\Gamma }`$ on $`L^2(\stackrel{~}{\mathrm{\Omega }},\stackrel{~}{F})`$. In fact it is easy to see that $`L^2(\stackrel{~}{\mathrm{\Omega }},\stackrel{~}{F})L^2\mathrm{\Gamma }L^2(U,\stackrel{~}{F})L^2\mathrm{\Gamma }L^2(\mathrm{\Omega },F)`$.
We have a unitary action of $`\mathrm{\Gamma }`$ on $`L^2\mathrm{\Gamma }`$ by left translations: $`\gamma l_\gamma `$ where $`l_\gamma f(x)=f(\gamma ^1x)`$ for $`x\mathrm{\Gamma }`$, $`fL^2\mathrm{\Gamma }`$. It induces an action on $`L^2(\stackrel{~}{\mathrm{\Omega }},\stackrel{~}{F})`$ by $`\gamma L_\gamma =l_\gamma \mathrm{Id}`$. Finally we denote by $`𝒟(.,.)`$ the various spaces of smooth compactly supported sections.
Let us consider a formally self-adjoint, strongly elliptic, positive differential operator $`P`$ on $`M`$ acting on sections of $`F`$. Denote by $`\stackrel{~}{P}`$ the $`\mathrm{\Gamma }`$–invariant differential operator which is its pull-back to $`\stackrel{~}{M}`$. From $`\stackrel{~}{P}`$ we construct the following operators: the Friedrichs extension in $`L^2(\stackrel{~}{\mathrm{\Omega }},\stackrel{~}{F})`$ of $`\stackrel{~}{P}`$ with domain $`𝒟(\stackrel{~}{\mathrm{\Omega }},\stackrel{~}{F})`$ and the Friedrichs extension in $`L^2(U,\stackrel{~}{F})`$ of $`\stackrel{~}{P}`$ with domain $`𝒟(U,\stackrel{~}{F})`$. From now on we denote these extensions by $`\stackrel{~}{P}`$ and $`P_0`$. They are closed self-adjoint positive operators. It is known that $`\stackrel{~}{P}`$ is also $`\mathrm{\Gamma }`$ invariant i.e. it commutes with all $`L_\gamma `$. This amounts of saying that $`E_\lambda `$ commutes with $`L_\gamma `$, $`\gamma \mathrm{\Gamma }`$, where $`(E_\lambda )_\lambda `$ is the spectral family of $`\stackrel{~}{P}`$. On the other hand the Rellich lemma tells that $`P_0`$ has compact resolvent and hence discrete spectrum. We will take the task of comparing the distribution of the two spectra. Namely since $`E_\lambda `$ is $`\mathrm{\Gamma }`$ invariant its image $`R(E_\lambda )`$ is a $`\mathrm{\Gamma }`$ invariant closed subspace of the free Hilbert $`\mathrm{\Gamma }`$–module $`L^2\mathrm{\Gamma }L^2(U,\stackrel{~}{F})L^2(\stackrel{~}{\mathrm{\Omega }},\stackrel{~}{F})`$. In general for any Hilbert space $``$ we call the Hilbert space $`L^2\mathrm{\Gamma }`$ a free Hilbert $`\mathrm{\Gamma }`$–module. The action of $`\mathrm{\Gamma }`$ is defined as above by $`\gamma L_\gamma =l_\gamma \mathrm{Id}`$. For $`\mathrm{\Gamma }`$ invariant closed spaces (called $`\mathrm{\Gamma }`$ modules) one can associate a positive, possibly infinite real number, called von Neumann or $`\mathrm{\Gamma }`$–dimension, denoted $`dim_\mathrm{\Gamma }`$. For notions involving the $`\mathrm{\Gamma }`$–dimension and linear algebra for $`\mathrm{\Gamma }`$–modules we refer the reader to \[At\], \[Sh\] and \[Ko\] (in the latter proofs from scratch are given). We give here the barest discussion of this score. Let us denote by $`𝒜_\mathrm{\Gamma }`$ the von Neumann algebra which consists of all bounded linear operators in $`L^2\mathrm{\Gamma }`$ which commute to the action of $`\mathrm{\Gamma }`$. To describe $`𝒜_\mathrm{\Gamma }`$ let us consider the von Neumann $`_\mathrm{\Gamma }`$ algebra of all bounded operators on $`L^2\mathrm{\Gamma }`$ which commute with all $`L_\gamma `$. It is generated by all right translations. If we consider the orthonormal basis $`(\delta _\gamma )_\gamma `$ in $`L^2\mathrm{\Gamma }`$ where $`\delta _\gamma `$ is the Dirac delta function at $`\gamma `$, then the matrix of any operator $`A_\mathrm{\Gamma }`$ has the property that all its diagonal elements are equal. Therefore we define a natural trace on $`_\mathrm{\Gamma }`$ as the diagonal element, that is, $`\mathrm{tr}_\mathrm{\Gamma }A=(A\delta _e,\delta _e)`$ where $`e`$ is the neutral element. Now $`𝒜_\mathrm{\Gamma }`$ is the tensor product of $`_\mathrm{\Gamma }`$ and the algebra $`()`$ of all bounded operators on $``$. If $`\mathrm{Tr}`$ is the usual trace on $`()`$ then we have a trace on $`𝒜_\mathrm{\Gamma }`$ by $`\mathrm{Tr}_\mathrm{\Gamma }=\mathrm{tr}_\mathrm{\Gamma }\mathrm{Tr}`$. For any $`\mathrm{\Gamma }`$ invariant space $`LL^2\mathrm{\Gamma }`$ i.e. for any $`\mathrm{\Gamma }`$–module, the projection $`P_L𝒜_\mathrm{\Gamma }`$ and we define $`dim_\mathrm{\Gamma }L=\mathrm{Tr}_\mathrm{\Gamma }P_L`$. Let us just remark for later use that if $`LL^2(\stackrel{~}{\mathrm{\Omega }},\stackrel{~}{F})`$ is a $`\mathrm{\Gamma }`$–module and $`f_i`$ is an orthonormal basis of $`L`$ then
$$dim_\mathrm{\Gamma }L=\underset{i}{}_U|f_i|^2.$$
$`(1.1)`$
We denote in the sequel $`N_\mathrm{\Gamma }(\lambda ,\stackrel{~}{P})=dim_\mathrm{\Gamma }R(E_\lambda )`$. Similary we consider the spectral distribution function (counting function) $`N(\lambda ,P_0)=dimR(E_\lambda ^0)`$ where $`E_\lambda ^0`$ is the spectral family of $`P_0`$; it equals the number of eigenvalues $`\lambda `$. We want to compare $`N_\mathrm{\Gamma }(\lambda ,\stackrel{~}{P})`$ and $`N(\lambda ,P_0)`$. For this purpose we use essentially the analysis of Shubin \[Sh\]. However there exist a difference in our method, namely we work at the beginning with model operator $`P_0`$ the operator $`\stackrel{~}{P}`$ itself with Dirichlet boundary conditions on $`U`$ whereas Shubin considers a direct sum of tangent operators to $`\stackrel{~}{P}`$. So we do not have to truncate from the outset the eigenfunctions of the model $`P_0`$. (See also Remark 1.3 in \[Sh\] and compare e.g. formulas (2.7), (2.8) or (3.6) from \[Sh\] with our corresponding formulas.) To begin with we need a variational principle.
###### Proposition 1.1(\cite{Sh})
Let $`\stackrel{~}{P}`$ be a $`\mathrm{\Gamma }`$ invariant self-adjoint positive operator on a free $`\mathrm{\Gamma }`$–module $`L^2\mathrm{\Gamma }`$ where $``$ is Hilbert space. Then
$$\begin{array}{ccc}\hfill N_\mathrm{\Gamma }(\lambda ,\stackrel{~}{P})=sup\{dim_\mathrm{\Gamma }LL\text{is a}\mathrm{\Gamma }\text{module}& \mathrm{Dom}(\stackrel{~}{Q}),\hfill & \\ \hfill \stackrel{~}{Q}(f,f)& \lambda f^2,fL\}(1.2)\hfill & \end{array}$$
where $`\stackrel{~}{Q}`$ is the quadratic form of $`\stackrel{~}{P}`$.
Recall that $`\stackrel{~}{Q}`$ is the closed symmetric quadratic given by $`\mathrm{Dom}(\stackrel{~}{Q})=\mathrm{Dom}(\stackrel{~}{P}^{1/2})`$, $`\stackrel{~}{Q}(u)=(\stackrel{~}{P}^{1/2}u,\stackrel{~}{P}^{1/2}u)`$. From the variational principle we deduce the following.
###### Proposition 1.2 (Estimate from below)
The counting functions for $`\stackrel{~}{P}`$ and $`P_0`$ satisfy the inequality
$$N_\mathrm{\Gamma }(\lambda ,\stackrel{~}{P})N(\lambda ,P_0),\lambda $$
$`(1.3)`$
###### Demonstration Proof
Let us denote by $`\lambda _0\lambda _1\mathrm{}`$ the spectrum of $`P_0`$. Let $`\{e_i\}_i`$ be an orthonormal basis of $`L^2(U,\stackrel{~}{F})`$ which consists of eigenfunctions of $`P_0`$ corresponding to the eigenvalues $`\{\lambda _i\}_i`$; if we let $`\stackrel{~}{e}_i=0`$ on $`\stackrel{~}{\mathrm{\Omega }}\overline{U}`$ and $`\stackrel{~}{e}_i=e_i`$ on $`U`$, $`\stackrel{~}{e}_i\mathrm{Dom}(\stackrel{~}{Q})`$, $`\{L_\gamma \stackrel{~}{e}_i\}_{i,\gamma }`$ is an orthonormal basis of $`L^2(\stackrel{~}{\mathrm{\Omega }},\stackrel{~}{F})`$ and $`\stackrel{~}{e}_{i,\gamma }=L_\gamma \stackrel{~}{e}_i\mathrm{Dom}(\stackrel{~}{Q})`$. We have $`\stackrel{~}{Q}(\stackrel{~}{e}_{i,\gamma },\stackrel{~}{e}_{i^{^{}},\gamma ^{^{}}})=\delta _{i,i^{^{}}}\delta _{\gamma ,\gamma ^{^{}}}\lambda _i`$. Let $`\mathrm{\Phi }_\lambda ^0`$ be the subspace spanned by $`\{e_i:\lambda _i\lambda \}`$ in $`L^2(U,\stackrel{~}{F})`$ and $`\mathrm{\Phi }_\lambda `$ the closed subspace spanned by $`\{\stackrel{~}{e}_{i,\gamma }:\lambda _i\lambda \}`$ in $`L^2(\stackrel{~}{\mathrm{\Omega }},\stackrel{~}{F})`$. Then by (1.1)
$$dim_\mathrm{\Gamma }\mathrm{\Phi }_\lambda =\underset{\lambda _i\lambda \gamma \mathrm{\Gamma }}{}_U|\stackrel{~}{e}_{i,\gamma }|^2=\underset{\lambda _i\lambda }{}e_i_U^2=dim\mathrm{\Phi }_\lambda ^0=N(\lambda ,P_0)$$
since $`\stackrel{~}{e}_{i,\gamma }|_U`$ vanishes unless $`\gamma `$ is the identity, and then it equals $`e_i`$ . If $`f`$ is a linear combination of $`\stackrel{~}{e}_{i,\gamma },\lambda _i\lambda `$, then $`\stackrel{~}{Q}(f,f)f^2`$ and, as $`\mathrm{Dom}(\stackrel{~}{Q})`$ is complete in the graph norm, we obtain that $`\mathrm{\Phi }_\lambda \mathrm{Dom}(\stackrel{~}{Q})`$ and $`\stackrel{~}{Q}(f,f)\lambda f^2`$, $`f\mathrm{\Phi }_\lambda `$. From the variational principle it follows that $`N_\mathrm{\Gamma }(\lambda ,\stackrel{~}{P})N(\lambda ,P_0)`$. ∎
The next step is an estimate from above of $`N_\mathrm{\Gamma }(\lambda ,\stackrel{~}{P})`$. Before let us say something about $`\mathrm{\Gamma }`$–morphisms. If $`L_1`$, $`L_2`$ are two $`\mathrm{\Gamma }`$–modules then an bounded liniar operator $`T:L_1L_2`$ is called a $`\mathrm{\Gamma }`$–morphism if it commutes with the action of $`\mathrm{\Gamma }`$. As for the usual dimension the following statements are true (see \[Ko\]). If $`T`$ is injective then $`dim_\mathrm{\Gamma }L_1dim_\mathrm{\Gamma }L_2`$ and if $`T`$ has dense image then $`dim_\mathrm{\Gamma }L_1dim_\mathrm{\Gamma }L_2`$. We denote by $`\mathrm{rank}_\mathrm{\Gamma }T=dim_\mathrm{\Gamma }\overline{R(T)}`$. For the following we refer to \[Sh\], Lemma 3.7.
###### Lemma 1.3
Let us consider the same setting as in the variational principle. Assume there is $`T:L^2(\stackrel{~}{\mathrm{\Omega }},\stackrel{~}{F})L^2(\stackrel{~}{\mathrm{\Omega }},\stackrel{~}{F})`$ a $`\mathrm{\Gamma }`$–morphism such that $`((\stackrel{~}{P}+T)f,f)\mu f^2`$, $`f\mathrm{Dom}(\stackrel{~}{P})`$ and $`\mathrm{rank}_\mathrm{\Gamma }Tp`$. Then
$$N_\mathrm{\Gamma }(\mu \epsilon ,\stackrel{~}{P})p,\epsilon >0.$$
$`(1.4)`$
In order to get an estimate from above we have to enlarge a little bit the fundamental domain $`U`$ and compare the counting function of $`\stackrel{~}{P}`$ to the counting function of the Friedrichs extension of $`\stackrel{~}{P}`$ restricted to compactly supported forms in the enlarged domain. For $`h>0`$, the enlarged domain is $`U_h=\{x\stackrel{~}{\mathrm{\Omega }}d(x,U)<h\}`$ where $`d`$ is the distance on $`\stackrel{~}{M}`$ associated to the Riemann metric on $`\stackrel{~}{M}`$. Then we take the tranlations $`U_{h,\gamma }:=\gamma U_h`$. Next we construct a partition of unity. Let $`\phi ^{(h)}C^{\mathrm{}}(\stackrel{~}{\mathrm{\Omega }})`$, $`\phi ^{(h)}0`$, $`\phi ^{(h)}=1`$ on $`\overline{U}`$ and $`\mathrm{supp}\phi ^{(h)}U_h`$, $`\phi _\gamma ^{(h)}=\phi ^{(h)}\gamma ^1`$. We define the function $`J_\gamma ^{(h)}C^{\mathrm{}}(\stackrel{~}{\mathrm{\Omega }})`$ by $`J_\gamma ^{(h)}=\phi _\gamma ^{(h)}\left(_\gamma (\phi _\gamma ^{(h)})^2\right)^{\frac{1}{2}}`$ so that $`_{\gamma \mathrm{\Gamma }}(J_\gamma ^{(h)})^2=1`$. If $`\stackrel{~}{P}`$ is of order $`2`$, which will be assumed throughout the section, then by \[Sh, Lemma 3.1\] (Shubin’s IMS localization formula, see \[CFKS\]) we know how to recover the operator $`\stackrel{~}{P}`$ from its localisations $`J_\gamma ^{(h)}\stackrel{~}{P}J_\gamma ^{(h)}`$:
$$\stackrel{~}{P}=\underset{\gamma \mathrm{\Gamma }}{}J_\gamma ^{(h)}\stackrel{~}{P}J_\gamma ^{(h)}\underset{\gamma \mathrm{\Gamma }}{}\sigma _0(\stackrel{~}{P})(dJ_\gamma ^{(h)})$$
$`(1.5)`$
where $`\sigma _0`$ is the principal symbol of $`\stackrel{~}{P}`$. In (1.5) $`J_\gamma ^{(h)}`$ are thought as multiplication operators on $`L^2(\stackrel{~}{\mathrm{\Omega }},\stackrel{~}{F})`$ – for which $`\mathrm{Dom}(\stackrel{~}{P})`$ is invariant – while $`_{\gamma \mathrm{\Gamma }}\sigma _0(\stackrel{~}{P})(dJ_\gamma ^{(h)})`$ is the multiplication by a bounded function. Since the derivative of $`J_\gamma ^{(h)}`$ is $`O(h^1)`$ and the order of $`\stackrel{~}{P}`$ is $`2`$ we see that the latter function is bounded by $`Ch^2`$ for some constant $`C>0`$ (here we use that the symbol is periodic and that $`\phi _\gamma ^{(h)}`$ are the translates of $`\phi ^{(h)}`$). Therefore the operatorial norm of the multiplication satisfies the same estimate and we deduce from (1.5) that
$$\stackrel{~}{P}\underset{\gamma \mathrm{\Gamma }}{}J_\gamma ^{(h)}\stackrel{~}{P}J_\gamma ^{(h)}\frac{C}{h^2}\mathrm{Id}$$
$`(1.6)`$
We consider now the operator $`\stackrel{~}{P}`$ with domain $`𝒟(U_h,\stackrel{~}{F})`$ and take its Friedrichs extension denoted $`P_0^{(h)}`$. We will compare $`N_\mathrm{\Gamma }(\lambda ,\stackrel{~}{P})`$ with the counting function of $`P_0^{(h)}`$. Let us fix $`\lambda `$. Denote by $`(E_\lambda ^{(h)})`$ the spectral family of $`P_0^{(h)}`$ and fix a positive constant $`M=M^{(\lambda )}`$ such that $`M\lambda inf\mathrm{spectrum}P_0^{(h)}`$ to the effect that $`P_0^{(h)}+ME_\lambda ^{(h)}\lambda \mathrm{Id}`$ . We define now a localisation of $`E_\lambda ^{(h)}`$ by taking the bounded operators $`G_\gamma ^{(h)}`$ on $`L^2(\stackrel{~}{\mathrm{\Omega }},\stackrel{~}{F})`$ given by $`G_\gamma ^{(h)}=J_\gamma ^{(h)}L_\gamma ME_\lambda ^{(h)}L_\gamma ^1J_\gamma ^{(h)}`$ and then summing over $`\mathrm{\Gamma }`$, $`G^{(h)}=_{\gamma \mathrm{\Gamma }}G_\gamma ^{(h)}`$. We have
$$\begin{array}{ccc}\hfill \stackrel{~}{P}+G^{(h)}& \underset{\gamma \mathrm{\Gamma }}{}\left(J_\gamma ^{(h)}\stackrel{~}{P}J_\gamma ^{(h)}+J_\gamma ^{(h)}L_\gamma ME_\lambda ^{(h)}L_\gamma ^1J_\gamma ^{(h)}\right)\frac{C}{h^2}\mathrm{Id}\hfill & \\ & =\underset{\gamma \mathrm{\Gamma }}{}J_\gamma ^{(h)}L_\gamma (H_0^{(h)}+ME_\lambda ^{(h)})L_\gamma ^1J_\gamma ^{(h)}\frac{C}{h^2}\mathrm{Id}(1.7)\hfill & \\ & \underset{\gamma \mathrm{\Gamma }}{}J_\gamma ^{(h)}L_\gamma \lambda L_\gamma ^1J_\gamma ^{(h)}\frac{C}{h^2}\mathrm{Id}=\left(\lambda \frac{C}{h^2}\right)\mathrm{Id}.\hfill & \end{array}$$
It is clear that $`G^{(h)}`$ will play the role of $`T`$ in Lemma 2.3. We must check one more hypothesis.
###### Claim 1.4
$$\mathrm{rank}_\mathrm{\Gamma }G^{(h)}N(\lambda ,P_0^{(h)})$$
$`(1.8)`$
###### Demonstration Proof
We start with the bounded operator $`\overline{G}^{(h)}`$ on $`L^2(U_s,\stackrel{~}{F})`$, given by $`\overline{G}^{(h)}=J_e^{(h)}ME_\lambda ^{(h)}J_e^{(h)}`$ . It is a finite rank operator, $`\mathrm{rank}\overline{G}^{(h)}\mathrm{rank}E_\lambda ^{(h)}=N(\lambda ,P_0^{(h)})`$. Next we consider the free $`\mathrm{\Gamma }`$–module $`L^2\mathrm{\Gamma }L^2(U_h,\stackrel{~}{F})`$ and the bounded $`\mathrm{\Gamma }`$–invariant operator $`\mathrm{Id}\overline{G}^{(h)}`$. Then $`R(\mathrm{Id}\overline{G}^{(h)})=L^2\mathrm{\Gamma }R(\overline{G}^{(h)})`$ so that $`\mathrm{rank}_\mathrm{\Gamma }\mathrm{Id}\overline{G}^{(h)}=\mathrm{rank}\overline{G}^{(h)}`$. We identify now the space $`L^2\mathrm{\Gamma }L^2(U_h,\stackrel{~}{F})`$ with $`_{\gamma \mathrm{\Gamma }}L^2(U_{h,\gamma },\stackrel{~}{F})`$ by the unitary transform $`K:_\gamma \delta _\gamma w_\gamma \left(L_\gamma w_\gamma \right)_\gamma `$ . Thus $`_{\gamma \mathrm{\Gamma }}L^2(U_{h,\gamma },\stackrel{~}{F})`$ is naturally a free $`\mathrm{\Gamma }`$–module for which $`K`$ is $`\mathrm{\Gamma }`$ invariant. We transport $`\mathrm{Id}\overline{G}^{(h)}`$ on $`_{\gamma \mathrm{\Gamma }}L^2(U_{h,\gamma },\stackrel{~}{F})`$ by $`K`$ and we think it as acting on this latter space. We construct then a restriction operator $`V:_{\gamma \mathrm{\Gamma }}L^2(U_{h,\gamma },\stackrel{~}{F})L^2(\stackrel{~}{\mathrm{\Omega }},\stackrel{~}{F})`$ , $`V\left((w_\gamma )_\gamma \right)=_{\gamma \mathrm{\Gamma }}w_\gamma `$ which is a surjective $`\mathrm{\Gamma }`$–morphism. We have also the $`\mathrm{\Gamma }`$–morphism $`I`$ from $`L^2(\stackrel{~}{\mathrm{\Omega }},\stackrel{~}{F})`$ to $`_{\gamma \mathrm{\Gamma }}L^2(U_{h,\gamma },\stackrel{~}{F})`$, $`I(u)=(u_{U_{h,\gamma }})_\gamma `$ which is obviously bounded. With our identifications we have $`G^{(h)}=V(\mathrm{Id}\overline{G}^{(h)})I`$ . As in the case of usual dimension $`\mathrm{rank}_\mathrm{\Gamma }V(\mathrm{Id}\overline{G}^{(h)})I\mathrm{rank}_\mathrm{\Gamma }(\mathrm{Id}\overline{G}^{(h)})`$ (see \[Sh\], Lemma 3.6).Therfore we conclude $`\mathrm{rank}_\mathrm{\Gamma }G^{(h)}\mathrm{rank}_\mathrm{\Gamma }(\mathrm{Id}\overline{G}^{(h)})=\mathrm{rank}\overline{G}^{(h)}N(\lambda ,P_0^{(h)})`$ . ∎
###### Proposition 1.5 (Estimate from above)
There is a constant $`C0`$ such that
$$N_\mathrm{\Gamma }(\lambda ,\stackrel{~}{P})N(\lambda +\frac{C}{h^2},P_0^{(h)})\lambda ,h>0$$
$`(1.9)`$
###### Demonstration Proof
The hypothesis of Lemma 2.3 are fulfilled for $`T=G^{(h)}`$, $`\mu =\lambda Ch^2`$ and $`p=N(\lambda ,P_0^{(h)})`$ as (1.7) and (1.8) show. Thus $`N_\mathrm{\Gamma }(\lambda \frac{C}{h^2}\epsilon ,\stackrel{~}{P})N(\lambda ,P_0^{(h)})`$, if $`\epsilon >0`$. Replacing $`\lambda `$ with $`\lambda +Ch^2+\epsilon `$, we obtain $`N_\mathrm{\Gamma }(\lambda ,\stackrel{~}{P})N(\lambda +\frac{C}{h^2}+\epsilon ,P_0^{(h)})`$. When $`\epsilon 0`$ the estimate (1.9) follows since the spectrum distribution function is right continuous by definition. ∎
The estimates from below and above for $`N_\mathrm{\Gamma }(\lambda ,\stackrel{~}{P})`$ enable us to study as a by–product the behaviour for $`\lambda \mathrm{}`$ to obtain the Weyl asymptotics for periodic operators (Shubin, see \[RSS\] and the references therein).
###### Corollary 1.6
If $`\stackrel{~}{P}`$ is a periodic, positive, second order elliptic operator as above then
$$\begin{array}{cc}\hfill \underset{\lambda \mathrm{}}{lim}\lambda ^{n/2}N_\mathrm{\Gamma }(\lambda ,\stackrel{~}{P})& =\underset{\lambda \mathrm{}}{lim}\lambda ^{n/2}N(\lambda ,P_0)\hfill \\ & =(2\pi )^n_U_{T_x^{}\stackrel{~}{M}}N(1,\sigma _0(\stackrel{~}{P})(x,\xi ))𝑑\xi 𝑑x\hfill \end{array}$$
where $`\sigma _0(\stackrel{~}{P})(x,\xi )\mathrm{Herm}(\stackrel{~}{F},\stackrel{~}{F})`$ is the principal symbol of $`\stackrel{~}{P}`$ and $`N(1,\sigma _0(\stackrel{~}{P})(x,\xi ))`$ is the counting function for the eigenvalues of this hermitian matrix.
###### Demonstration Proof
First let us remark that the last equality is the classical Weyl type formula as established by Carleman, Gårding and others, see \[RSS\], p.72. It is obvious that $`lim\; inf\lambda ^{n/2}N(\lambda ,P_0)lim\; inf\lambda ^{n/2}N_\mathrm{\Gamma }(\lambda ,\stackrel{~}{P})`$ by the estimate from below. On the other hand the estimate from above gives
$$\begin{array}{cc}\hfill lim\; sup\lambda ^{n/2}N_\mathrm{\Gamma }(\lambda ,\stackrel{~}{P})lim\; sup\left(1+\frac{C}{\lambda h^2}\right)^{n/2}\left(\lambda +\frac{C}{h^2}\right)^{n/2}N(\lambda +\frac{C}{h^2},P_0^{(h)})& \\ \hfill lim\; sup\mu ^{n/2}N(\mu ,P_0^{(h)})=(2\pi )^n_{U_h}_{T_x^{}\stackrel{~}{M}}N(1,\sigma _0(\stackrel{~}{P})(x,\xi ))𝑑\xi 𝑑x& \end{array}$$
for a fixed small $`h`$. We make $`h0`$ and obtain the desired formula. ∎
We are going to apply the above results to the semi-classical asymptotics as $`k\mathrm{}`$ of the spectral distribution function of the laplacian $`k^1\stackrel{~}{\mathrm{\Delta }}^{\prime \prime }`$ on $`\stackrel{~}{M}`$. Let $`G`$ be a hermitian holomorphic bundle on $`M`$ and $`\stackrel{~}{G}=p^{}G`$ its pull-back. We define $`𝒟^{(0,q)}(.,.)`$ to be the space of smooth compactly supported $`(0,q)`$ forms. Let $`\overline{}:𝒟^{0,q}(\stackrel{~}{M},\stackrel{~}{G})𝒟^{0,q+1}(\stackrel{~}{M},\stackrel{~}{G})`$ be the Cauchy–Riemann operator and $`\vartheta :𝒟^{0,q+1}(\stackrel{~}{M},\stackrel{~}{G})𝒟^{0,q}(\stackrel{~}{M},\stackrel{~}{G})`$ the formal adjoint of $`\overline{}`$ with respect to the given hermitian metrics on $`\stackrel{~}{M}`$, $`\stackrel{~}{G}`$. Then $`\stackrel{~}{\mathrm{\Delta }}^{\prime \prime }=\overline{}\vartheta +\vartheta \overline{}`$ is a formally self-adjoint, strongly elliptic, positive and $`\mathrm{\Gamma }`$–invariant differential operator.
We take $`\stackrel{~}{E}`$ and $`\stackrel{~}{G}`$ two $`\mathrm{\Gamma }`$ invariant holomorphic bundles. Let us form the Laplace–Beltrami operator $`\stackrel{~}{\mathrm{\Delta }}_k^{\prime \prime }`$ on $`(0,q)`$ forms with values in $`\stackrel{~}{E}^k\stackrel{~}{G}`$. Thus we will consider the $`\mathrm{\Gamma }`$ invariant hermitian bundle $`\stackrel{~}{F}=\mathrm{\Lambda }^{(0,q)}T^{}\stackrel{~}{M}\stackrel{~}{E}^k\stackrel{~}{G}`$ and apply the previous results for $`\stackrel{~}{P}=k^1\stackrel{~}{\mathrm{\Delta }}_k^{\prime \prime }_{\stackrel{~}{\mathrm{\Omega }}}`$ where the index $`\stackrel{~}{\mathrm{\Omega }}`$ emphasises that the Friedrichs extension gives the operator of the Dirichlet problem on $`\stackrel{~}{\mathrm{\Omega }}`$. Now we have to make a good choice of the parameter $`h`$. We take $`h=k^{\frac{1}{4}}`$ so that the derivative of the cutting off function $`J_\gamma ^{(h)}`$ is just $`O(k^{\frac{1}{4}})`$. Then $`\sigma _0(k^1\stackrel{~}{\mathrm{\Delta }}_k^{\prime \prime })(dJ_\gamma ^{(h)})=k^1|\overline{}J_\gamma ^{(h)}|^2=O(k^{\frac{1}{2}})`$. Therefore formula (1.6) becomes $`\frac{1}{k}\stackrel{~}{\mathrm{\Delta }}_k^{\prime \prime }_{\stackrel{~}{\mathrm{\Omega }}}_{\gamma \mathrm{\Gamma }}J_\gamma ^{(h)}\frac{1}{k}\stackrel{~}{\mathrm{\Delta }}_k^{\prime \prime }_{\stackrel{~}{\mathrm{\Omega }}}J_\gamma ^{(h)}\frac{C}{\sqrt{k}}\mathrm{Id}.`$ We have thus proved the following semi–classical estimate for laplacian.
###### Proposition 1.7
There exists a constant $`C>0`$ such that for $`\lambda `$ and $`k>0`$ we have
$$N(\lambda ,\frac{1}{k}\stackrel{~}{\mathrm{\Delta }}_k^{\prime \prime }_U)N_\mathrm{\Gamma }(\lambda ,\frac{1}{k}\stackrel{~}{\mathrm{\Delta }}_k^{\prime \prime }_{\stackrel{~}{\mathrm{\Omega }}})N(\lambda +\frac{C}{\sqrt{k}},\frac{1}{k}\stackrel{~}{\mathrm{\Delta }}_k^{\prime \prime }_{U_{k^{1/4}}})$$
$`(1.10)`$
Demailly has determined the distribution of spectrum for the Dirichlet problem for $`\stackrel{~}{\mathrm{\Delta }}_k^{\prime \prime }`$ in \[De1\], Theorem 3.14. For this purpose he introduces (\[De1\],(1.5)) the function $`\nu _E:\stackrel{~}{M}\times `$ depending on the curvature of $`\stackrel{~}{E}`$ and then considers the function $`\overline{\nu }_E(x,\lambda )=lim_{\epsilon 0}\nu _E(x,\lambda +\epsilon )`$. The function $`\overline{\nu }_E(x,\lambda )`$ is right continuous in $`\lambda `$ and bounded above on compacts of $`\stackrel{~}{M}`$. Denote by $`\alpha _1(x),\mathrm{},\alpha _n(x)`$ the eigenvalues of of the curvature form $`i𝕔(\stackrel{~}{E})(x)`$ with respect to the metric on $`\stackrel{~}{M}`$. We also denote for a multiindex $`J\{1,\mathrm{},n\}`$, $`\alpha _J=_{jJ}\alpha _j`$ and $`C(J)=\{1,\mathrm{},n\}J`$ . For $`VM`$ we introduce
$$I^q(V,\mu )=\underset{J=q}{}_V\overline{\nu }_E(2\mu +\alpha _{C(J)}\alpha _J)𝑑\sigma $$
###### Proposition 1.8 (Demailly)
Assume that $`V`$ has measure zero and that the laplacian acts on $`(0,q)`$ forms. Then $`lim\; sup_kk^nN(\lambda ,\frac{1}{k}\mathrm{\Delta }_k^{\prime \prime }_V)I^q(V,\lambda )`$ Moreover there exists an at most countable set $`𝒩`$ such that for $`\lambda 𝒩`$ the limit of the left–hand side expression exists and we have equality.
We return now to the case of a covering manifold and apply Demailly’s formula in (1.10). Let us fix $`\epsilon >0`$. For sufficiently large $`k`$ we have $`U_{k^{\frac{1}{4}}}U_\epsilon `$ so the fact that the counting function is increasing and the variational principle yield $`N(\lambda +\frac{C}{\sqrt{k}},\frac{1}{k}\stackrel{~}{\mathrm{\Delta }}^{\prime \prime }_{U_{k^{1/4}}})N(\lambda +\epsilon ,\frac{1}{k}\stackrel{~}{\mathrm{\Delta }}^{\prime \prime }_{U_{k^{1/4}}})N(\lambda +\epsilon ,\frac{1}{k}\stackrel{~}{\mathrm{\Delta }}^{\prime \prime }_{U_\epsilon })`$. Hence by (1.10) and Proposition 1.8 ($`U_\epsilon `$ is negligible for small $`\epsilon `$),
$$\underset{k}{lim\; sup}k^nN_\mathrm{\Gamma }(\lambda ,\frac{1}{k}\stackrel{~}{\mathrm{\Delta }}_k^{\prime \prime }_{\stackrel{~}{\mathrm{\Omega }}})I^q(U_\epsilon ,\lambda +\epsilon ).$$
The use of dominated convergence to make $`\epsilon 0`$ in the last integral yield the asymptotic formula for the laplacian on a covering manifold.
###### Theorem 1.9
The spectral distribution function of $`\frac{1}{k}\stackrel{~}{\mathrm{\Delta }}_k^{\prime \prime }_{\stackrel{~}{\mathrm{\Omega }}}`$ on $`L_{0,q}^2(\stackrel{~}{\mathrm{\Omega }},\stackrel{~}{E}^k\stackrel{~}{G})`$ with Dirichlet boundary conditions satisfies
$$\underset{k}{lim\; sup}k^nN_\mathrm{\Gamma }(\lambda ,\frac{1}{k}\stackrel{~}{\mathrm{\Delta }}_k^{\prime \prime }_{\stackrel{~}{\mathrm{\Omega }}})I^q(U,\lambda ).$$
$`(1.11)`$
Moreover, there exists an at most countable set $`𝒩`$ such that for $`\lambda 𝒩`$ the limit exits and we have equality in rm (1.11).
## §2 Geometric situations
In this section we apply the results from the previous section to the study of the $`L^2`$ cohomology of coverings of complex manifolds satisfying certain curvature conditions. If $`M`$ is a complete Kähler manifold and $`E`$ a positive line bundle on $`M`$ the $`L^2`$ estimates of Andreotti–Vesentini–Hörmander allow to find a lot of sections of $`\stackrel{~}{E}`$ on a covering $`\stackrel{~}{M}`$ (see e.g. \[NR\]). We prove here the following.
###### Theorem 2.1
Let $`(M,\omega )`$ be an $`n`$–dimensional complete hermitian manifold such that the torsion of $`\omega `$ is bounded and let $`(E,h)`$ be a holomorphic hermitian line bundle. Let $`KM`$ and a constant $`C_0>0`$ such that $`ı𝚌(E,h)C_0\omega `$ on $`MK`$. Let $`p:\stackrel{~}{M}M`$ be a Galois covering with group $`\mathrm{\Gamma }`$ and $`\stackrel{~}{E}=p^{}E`$ and let $`\mathrm{\Omega }`$ be any open set with smooth boundary and $`K\mathrm{\Omega }M`$. Then
$$dim_\mathrm{\Gamma }H_{(2)}^{n,0}(\stackrel{~}{M},\stackrel{~}{E}^k)\frac{k^n}{n!}_{\mathrm{\Omega }(1,h)}\left(\frac{ı}{2\pi }𝚌(E,h)\right)^n+o(k^n),k>>0,$$
where $`H_{(2)}^{n,0}(\stackrel{~}{M},\stackrel{~}{E}^k)`$ is the space of $`(n,0)`$–forms with values in $`\stackrel{~}{E}^k`$ which are $`L^2`$ with respect to any metric on $`\stackrel{~}{M}`$ and the pullback of $`h`$ and $`\mathrm{\Omega }(1,h)`$ is the subset of $`\mathrm{\Omega }`$ where $`ı𝚌(E,h)`$ is non–degenerate and has at most one negative eigenvalue.
###### Demonstration Proof
We endow $`\stackrel{~}{M}`$ with the metric $`\stackrel{~}{\omega }=p^{}\omega `$ and $`\stackrel{~}{E}`$ with $`\stackrel{~}{h}=p^{}h`$. All the norms, Laplace–Beltrami operators, spaces of harmonic forms and $`L^2`$–cohomology groups are with respect to $`\stackrel{~}{\omega }`$ and $`\stackrel{~}{h}`$. In particular the operators $`\overline{}`$ and Lapalce–Beltrami are $`\mathrm{\Gamma }`$–invariant. It is standard to see that $`\stackrel{~}{\omega }`$ is also complete. To justify this let us first take a compact set $`KX`$ and consider $`\stackrel{~}{K}=p^1K`$. The metric $`\stackrel{~}{\omega }_\epsilon `$ is complete on $`\stackrel{~}{K}`$ in the following sense. There exist functions $`\phi _\epsilon 𝒞^{\mathrm{}}(\stackrel{~}{K})`$ with values in $`[0,1]`$ such that $`\mathrm{supp}\phi _\epsilon `$ is compact in $`\stackrel{~}{K}`$, the sets $`\{z\stackrel{~}{K}:\phi _\epsilon (z)=1\}`$ form an exhaustion of $`\stackrel{~}{K}`$ and $`sup|d\phi _\epsilon |=O(\epsilon )`$ as $`\epsilon 0`$. This is seen as usual by observing that the balls are relatively compact in $`\stackrel{~}{K}`$ and then taking cut–off functions. Since $`M`$ is complete there exist an exhaustion $`K_\nu `$ with compacts and functions $`\psi _\nu 𝒞^{\mathrm{}}(M)`$ with values in $`[0,1]`$ and $`\mathrm{supp}\psi _\nu K_{\nu +1}`$ such that $`K_\nu \{zM:\psi _\nu (z)=1\}`$ and $`sup|d\psi _\nu |2^\nu `$. Let us choose now a point $`z_0\stackrel{~}{K}_0`$ and fix fundamental domains $`U_\nu `$ for the action of $`\mathrm{\Gamma }`$ on $`\stackrel{~}{K}_\nu `$ such that $`z_0U_\nu `$. We also choose an exhaustion by finite sets $`I_0I_1\mathrm{}I_\nu \mathrm{}\mathrm{\Gamma }`$ of $`\mathrm{\Gamma }`$. Indeed, since $`\stackrel{~}{M}`$ is paracompact $`\mathrm{\Gamma }`$ is countable. For each $`\nu `$ let us take $`\phi _\nu 𝒞^{\mathrm{}}(\stackrel{~}{K}_{\nu +1})`$ such that $`\phi _\nu =1`$ on $`\{\gamma U_{\nu +1}:\gamma I_{\nu +1}\}`$ and $`sup|d\phi _\nu |2^\nu `$. We consider also the function $`\stackrel{~}{\psi }_\nu =\psi _\nu p`$. Then the functions $`\stackrel{~}{\psi }_\nu \phi _\nu `$ have compact support in $`\stackrel{~}{M}`$, the sets where they equal $`1`$ exhaust $`\stackrel{~}{M}`$ and their derivative is $`O(2^{\nu 1})`$, which proves that $`\stackrel{~}{M}`$ is complete. We remark here that $`𝒰=_\nu U_\nu `$ is a fundamental domain for the action of $`\mathrm{\Gamma }`$ on $`\stackrel{~}{M}`$ and that if $`\stackrel{~}{G}`$ is a $`\mathrm{\Gamma }`$–invariant bundle on $`\stackrel{~}{M}`$ then $`L^2(\stackrel{~}{M},\stackrel{~}{G})`$ is a free $`\mathrm{\Gamma }`$–module.
We take $`\mathrm{\Omega }`$ as in the hypothesis and let $`U`$ be a fundamental domain of $`\stackrel{~}{\mathrm{\Omega }}`$ as in §1. Since $`p`$ is locally biholomorphic we see that $`ı𝕔(\stackrel{~}{E},\stackrel{~}{h})C_0\stackrel{~}{\omega }`$ on $`\stackrel{~}{M}\stackrel{~}{K}`$. Let $`u`$ be a smooth $`(n,1)`$ form on $`\stackrel{~}{M}`$ with values in $`\stackrel{~}{E}^k`$ and compactly supported outside $`\stackrel{~}{K}`$. We apply now the Bochner–Kodaira–Nakano formula for $`u`$:
$$3(\stackrel{~}{\mathrm{\Delta }}_k^{\prime \prime }u,u)2([ı𝚌(\stackrel{~}{E}^k),\stackrel{~}{\mathrm{\Lambda }}]u,u)\left(\tau u^2+\overline{\tau }u^2+\tau ^{}u^2+\overline{\tau }^{}u^2\right),$$
where $`\mathrm{\Lambda }`$ is the operator of taking the interior product with $`\stackrel{~}{\omega }`$ and the $`\tau `$’s are the torsion operators of the metric $`\stackrel{~}{\omega }`$. More precisely $`\tau =[\mathrm{\Lambda },\stackrel{~}{\omega }]`$. Therefore there exists a constant $`C_1>0`$ (depending just on the metric $`\omega `$) such that
$$3(\stackrel{~}{\mathrm{\Delta }}_k^{\prime \prime }u,u)2C_0ku^2C_1u^2,$$
and hence
$$(\stackrel{~}{\mathrm{\Delta }}_k^{\prime \prime }u,u)\frac{C_0k}{2}u^2,k\frac{C_1}{2C_0}.$$
$`(2.1)`$
Indeed, by hypothesis the torsion operators are pointwise bounded. Moreover $`([ı𝚌(\stackrel{~}{E}^k,\stackrel{~}{h}^k),\mathrm{\Lambda }]u,u)k\alpha _1|u|^2`$ where $`\alpha _1\mathrm{}\alpha _n`$ are the eigenvalues of $`ı𝚌(\stackrel{~}{E},\stackrel{~}{h})`$ with respect to $`\stackrel{~}{\omega }`$.
Let $`\rho 𝒞^{\mathrm{}}(M)`$ such that $`\rho =0`$ on $`L`$ and $`\rho =1`$ on $`M\mathrm{\Omega }`$, where $`L`$ is a neighbourhood of $`K`$ in $`\mathrm{\Omega }`$. We put $`\stackrel{~}{\rho }=\rho p`$. Let $`u𝒟^{n,1}(\stackrel{~}{M},\stackrel{~}{E}^k)`$, so that $`\stackrel{~}{\rho }u`$ has support outside $`\stackrel{~}{K}`$. We use now the elementary estimate:
$$(\stackrel{~}{\mathrm{\Delta }}_k^{\prime \prime }(\stackrel{~}{\rho }u),\stackrel{~}{\rho }u)\frac{3}{2}(\stackrel{~}{\mathrm{\Delta }}_k^{\prime \prime }u,u)+6sup|d\stackrel{~}{\rho }|^2u^2.$$
$`(2.2)`$
Obviously $`C_2=6sup|d\stackrel{~}{\rho }|^2<\mathrm{}`$. Estimates (2.1) and (2.2) yield
$$u^2\frac{12}{C_0k}(\stackrel{~}{\mathrm{\Delta }}_k^{\prime \prime }u,u)+4_{\stackrel{~}{\mathrm{\Omega }}}\left|(1\stackrel{~}{\rho })u\right|^2,k\frac{\mathrm{max}\{C_1,16C_2\}}{2C_0}$$
$`(2.3)`$
for any compactly supported $`u`$. Since the metric $`\stackrel{~}{\omega }`$ is complete the density lemma of Andreotti and Vesentini \[AV\] shows that $`\stackrel{~}{\mathrm{\Delta }}_k^{\prime \prime }`$ is essentially self–adjoint. Thus (2.3) is true for any $`u`$ in the domain of the quadratic form $`\stackrel{~}{Q}_k`$ of the self–adjoint extension of $`k^1\stackrel{~}{\mathrm{\Delta }}_k^{\prime \prime }`$. From relation (2.3) we infer that the spectral spaces corresponding to the lower part of the spectrum of $`k^1\stackrel{~}{\mathrm{\Delta }}_k^{\prime \prime }`$ on $`(n,1)`$–forms can be injected into the spectral spaces of the $`\mathrm{\Gamma }`$–invariant operator $`k^1\stackrel{~}{\mathrm{\Delta }}_k^{\prime \prime }_{\stackrel{~}{\mathrm{\Omega }}}`$ which correspond to the Dirichlet problem on $`\stackrel{~}{\mathrm{\Omega }}`$ for $`k^1\stackrel{~}{\mathrm{\Delta }}_k^{\prime \prime }`$. The latter operator was studied in §1. This idea appears in Witten’s proof (see Henniart \[He\]) and in \[Bou\] in the context of $`q`$–convex manifolds in the sense of Andreotti–Grauert. We claim that for $`\lambda <C_0/24`$,
$$L_k^1(\lambda )L_{k,\stackrel{~}{\mathrm{\Omega }}}^1(12\lambda +C_3k^1),uE_{12\lambda +C_3k^1}(k^1\stackrel{~}{\mathrm{\Delta }}_k^{\prime \prime }_{\stackrel{~}{\mathrm{\Omega }}})(1\stackrel{~}{\rho })u,$$
$`(2.4)`$
is an injective $`\mathrm{\Gamma }`$–morphism, where $`L_k^1(\lambda )=\mathrm{Range}\left(E_\lambda (k^1\stackrel{~}{\mathrm{\Delta }}_k^{\prime \prime }_{\stackrel{~}{\mathrm{\Omega }}})\right)`$ is the spectral space of $`k^1\stackrel{~}{\mathrm{\Delta }}_k^{\prime \prime }`$ on $`(n,1)`$–forms, $`L_{k,\stackrel{~}{\mathrm{\Omega }}}^1(\mu )=\mathrm{Range}E_\mu (k^1\stackrel{~}{\mathrm{\Delta }}_k^{\prime \prime }_{\stackrel{~}{\mathrm{\Omega }}})`$, the spectral spaces of $`k^1\stackrel{~}{\mathrm{\Delta }}_k^{\prime \prime }_{\stackrel{~}{\mathrm{\Omega }}}`$ and $`C_3=8C_2`$. To prove the claim let us remark that the map (2.4) is the restriction of an operator on $`L_{0,1}^2(\stackrel{~}{M},\stackrel{~}{E}^kK_{\stackrel{~}{M}})`$ of the same form; this is continuous and $`\mathrm{\Gamma }`$–invariant being a composition of a multiplication with a bounded $`\mathrm{\Gamma }`$–invariant function and a $`\mathrm{\Gamma }`$–invariant projection. To prove the injectivity we choose $`uL_k^1(\lambda )`$, $`\lambda <C_0/24`$ to the effect that $`\stackrel{~}{Q}_k(u)\lambda u^2(C_0/24)u^2`$. Plugging this relation in (2.3) we get
$$u^28_{\stackrel{~}{\mathrm{\Omega }}}\left|(1\stackrel{~}{\rho })u\right|^2,uL_k^1(\lambda ),\lambda <C_0/24.$$
$`(2.5)`$
Let us denote by $`\stackrel{~}{Q}_{k,\stackrel{~}{\mathrm{\Omega }}}`$ the quadratic form of $`k^1\stackrel{~}{\mathrm{\Delta }}_k^{\prime \prime }_{\stackrel{~}{\mathrm{\Omega }}}`$. Then by (2.2) and (2.5), $`\stackrel{~}{Q}_{k,\stackrel{~}{\mathrm{\Omega }}}\left((1\stackrel{~}{\rho })u\right)\frac{3}{2}\stackrel{~}{Q}_k(u)+\frac{C_2}{k}u^2\left(12\lambda +\frac{8C_2}{k}\right)_{\stackrel{~}{\mathrm{\Omega }}}\left|(1\stackrel{~}{\rho })u\right|^2`$ which shows that if $`E(12\lambda +C_3k^1,k^1\stackrel{~}{\mathrm{\Delta }}_k^{\prime \prime }_{\stackrel{~}{\mathrm{\Omega }}})(1\stackrel{~}{\rho })u=0`$ then $`(1\stackrel{~}{\rho })u=0`$ so that $`u=0`$ by (2.5). Therefore (2.4) is injective and hence
$$N_\mathrm{\Gamma }^1(\lambda ,\frac{1}{k}\stackrel{~}{\mathrm{\Delta }}_k^{\prime \prime })N_\mathrm{\Gamma }^1(12\lambda +\frac{C_3}{k},\stackrel{~}{\mathrm{\Delta }}_k^{\prime \prime }_{\stackrel{~}{\mathrm{\Omega }}}),\lambda <(C_0/24),$$
$`(2.6)`$
and thus the spectral spaces $`L_k^1(\lambda )`$, $`\lambda <C_0/24`$, are of finite $`\mathrm{\Gamma }`$–dimension.
Now we can apply Theorem 1.9 for $`k^1\stackrel{~}{\mathrm{\Delta }}_k^{\prime \prime }_{\stackrel{~}{\mathrm{\Omega }}}`$ on $`\stackrel{~}{\mathrm{\Omega }}`$ (with $`\stackrel{~}{G}=K_{\stackrel{~}{M}}`$). By the variational principle we have that $`N_\mathrm{\Gamma }^0(\lambda ,\frac{1}{k}\stackrel{~}{\mathrm{\Delta }}_k^{\prime \prime })N_\mathrm{\Gamma }^0(\lambda ,\frac{1}{k}\stackrel{~}{\mathrm{\Delta }}_k^{\prime \prime }_{\stackrel{~}{\mathrm{\Omega }}})`$ and by Theorem 1.9 for $`q=0`$
$$\underset{k}{lim\; inf}k^nN_\mathrm{\Gamma }^0(\lambda ,\frac{1}{k}\stackrel{~}{\mathrm{\Delta }}_k^{\prime \prime })I^0(U,\lambda ),\lambda <C_0/24,\lambda 𝒩$$
$`(2.7)`$
We find now an upper bound. Fix an arbitrary $`\delta >0`$. For $`k>C_3/\delta `$ we have $`N^1(\lambda ,k^1\stackrel{~}{\mathrm{\Delta }}_k^{\prime \prime })N_\mathrm{\Gamma }^1(12\lambda +C_3k^1\stackrel{~}{\mathrm{\Delta }}_k^{\prime \prime }_{\stackrel{~}{\mathrm{\Omega }}})N_\mathrm{\Gamma }^1(12\lambda +\delta ,\stackrel{~}{\mathrm{\Delta }}_k^{\prime \prime }_{\stackrel{~}{\mathrm{\Omega }}})`$ hence by (1.11) $`lim\; sup_kk^nN_\mathrm{\Gamma }^1(\lambda ,\frac{1}{k}\stackrel{~}{\mathrm{\Delta }}_k^{\prime \prime })I^1(U,12\lambda +\delta )`$. We can let $`\delta 0`$ so that
$$\underset{k}{lim\; sup}k^nN_\mathrm{\Gamma }^1(\lambda ,\frac{1}{k}\stackrel{~}{\mathrm{\Delta }}_k^{\prime \prime })I^1(U,12\lambda ),\lambda <C_0/24.$$
$`(2.8)`$
We consider the group $`H_{(2)}^{n,0}(\stackrel{~}{M},\stackrel{~}{E}^k)=\{uL_{n,0}^2(\stackrel{~}{M},\stackrel{~}{E}^k,\stackrel{~}{\omega },\stackrel{~}{h}):\overline{}u=0\}`$ which is a $`\mathrm{\Gamma }`$–module and we find a lower bound for its $`\mathrm{\Gamma }`$–dimension. We know that the $`L^2`$ norm doesn’t actually depend on the metric on $`\stackrel{~}{M}`$. We consider also the operator $`\stackrel{~}{\mathrm{\Delta }}_k^{\prime \prime }`$ defined on $`L_{n,0}^2(\stackrel{~}{M},\stackrel{~}{E}^k)`$ and denote by $`L_k^0(\lambda )`$ its spectral spaces. Since $`\stackrel{~}{\mathrm{\Delta }}_k^{\prime \prime }`$ commutes with $`\overline{}`$ it follows that the spectral projections of $`\stackrel{~}{\mathrm{\Delta }}_k^{\prime \prime }`$ commute with $`\overline{}`$ too, showing thus $`\overline{}L_k^0(\lambda )L_k^1(\lambda )`$ and therefore we have the $`\mathrm{\Gamma }`$–morphism $`L_k^0(\lambda )\mathrm{@}>\overline{}_\lambda >>L_k^1(\lambda )`$ where $`\overline{}_\lambda `$ denotes the restriction of $`\overline{}`$ (by the definition of $`L_k^0(\lambda )`$, $`\overline{}_\lambda `$ is bounded by $`k\lambda `$). Since for any $`\mathrm{\Gamma }`$–morphism $`A`$ we have $`dim_\mathrm{\Gamma }\overline{R(A)}=dim_\mathrm{\Gamma }\mathrm{ker}(A)^{}`$ we see that $`dim_\mathrm{\Gamma }\mathrm{ker}\overline{}_\lambda +dim_\mathrm{\Gamma }\overline{R(\overline{}_\lambda )}=dim_\mathrm{\Gamma }L_k^0(\lambda )`$. Moreover $`dim_\mathrm{\Gamma }\overline{R(\overline{}_\lambda )}dim_\mathrm{\Gamma }L_k^1(\lambda )`$ and they are finite. Therefore by (2.7) and (2.8), $`dim_\mathrm{\Gamma }H_{(2)}^{n,0}(M,\stackrel{~}{E}^k)dim_\mathrm{\Gamma }\mathrm{ker}\overline{}_\lambda k^n\left[I^0(U,2\lambda )I^1(U,12\lambda )\right]`$ for $`\lambda <C_0/24`$ and $`\lambda 𝒩`$. We can now let $`\lambda `$ go to zero through these values. The limits $`I^0(U,0)`$ and $`I^1(U,0)`$ are calculated in \[De1\] and if we identify the fundamental domain $`U`$ with $`\mathrm{\Omega }`$ the result is exactly the integral from the conclusion. ∎
To state the following result let us remind that by the definition of Andreotti and Grauert \[AG\] a manifold is called $`1`$– concave if there exists a smooth function $`\phi :X(a,b]`$ where $`a\{\mathrm{}\}`$, $`b`$, such that $`X_c:=\{\phi >c\}X`$ for all $`c(a,b]`$ and $`\phi `$ is strictly plurisubharmonic outside a compact set. Let $`E`$ be a holomorphic line bundle on $`X`$. In \[Oh\], \[Ma\] one constructs a function $`\chi :(\mathrm{},0)`$ such that $`_1^0\chi (t)^{1/2}𝑑t=\mathrm{}`$, $`\chi ^{}(t)^24\chi (t)^3`$ , $`\chi (t)4`$ and a hermitian metric $`\omega `$ which equals $`\frac{1}{3}\overline{}\phi `$ near $`bX_c`$. For convenience we denote $`\psi =c\phi `$. We define $`\omega _0=\omega +\chi (\psi )\phi \overline{}\phi `$, a complete metric on $`X_c`$ and a hermitian metric $`h_0=h\mathrm{exp}(A_{inf\psi }^\psi \chi (t)𝑑t)`$ on $`E`$ over $`X_c`$.
###### Theorem 2.2
Let $`X`$ be a $`1`$–concave manifold of dimension $`n3`$ and let $`X_c`$ be a sublevel set such that the exhaustion function $`\phi `$ is strictly plurisubharmonic near $`bX_c`$. Let $`p:\stackrel{~}{X}_cX_c`$ be a Galois covering of group $`\mathrm{\Gamma }`$. Assume that $`\stackrel{~}{X}_c`$ and $`\stackrel{~}{E}`$ are endowed with the lifts of the metrics $`\omega _0`$ and $`h_0`$. Then
$$dim_\mathrm{\Gamma }H_{(2)}^0(\stackrel{~}{X}_c,\stackrel{~}{E}^k)\frac{k^n}{n!}_{\mathrm{\Omega }(1,h_0)}\left(\frac{ı}{2\pi }𝚌(E,h_0)\right)^n+o(k^n),k>>0.$$
$`2.9`$
for any sufficiently large open set $`\mathrm{\Omega }X_c`$.
###### Demonstration Proof
The metrics $`\omega _0`$ and $`h_0`$ satisfy the following conditions:
(i) Denoting by $`\gamma _i`$ the eigenvalues of $`ı\chi (\psi )\overline{}\psi +ı\chi ^{}(\psi )\psi \overline{}\psi `$ with respect to $`\omega _0`$ we have $`\gamma _1\mathrm{}\gamma _{n1}2\chi (\psi )`$ and $`\gamma _n\chi (\psi )`$ so that $`\gamma _n+\mathrm{}+\gamma _2(52n)\chi (\psi )\chi (\psi )`$ for $`n3`$ outside a compact set $`K:=X_eX_c`$.
(ii) The torsion operators of the metric $`\omega _0`$ are pointwise bounded by $`C_2\chi (\phi )^{1/2}`$ outside $`K`$.
(iii) The eigenvalues of $`ı𝚌(E,h_0)`$ with respect to $`\omega _0`$ are bounded above on $`X_c`$ by $`C_1>0`$.
Let us take the lifts $`\stackrel{~}{\omega }_0`$, $`h_0`$ and $`\stackrel{~}{\psi }=c\phi p`$. It is easy to see that properties (i), (ii) and (iii) are still valid for $`\stackrel{~}{\omega }_0`$ and $`\stackrel{~}{h}_0`$ and $`\stackrel{~}{\psi }`$ on $`\stackrel{~}{X}_c\stackrel{~}{K}`$. For $`u𝒟^{(0,1)}(\stackrel{~}{X}_c\stackrel{~}{K},\stackrel{~}{E}^k)`$ we apply the Bochner–Kodaira–Nakano inequality and take into account the formula $`([ı𝚌(\stackrel{~}{E}^k,\stackrel{~}{h}^k),\mathrm{\Lambda }]u,u)k(\alpha _n+\mathrm{}+\alpha _2)|u|^2`$. Then
$$3(\stackrel{~}{\mathrm{\Delta }}_k^{\prime \prime }u,u)\left(knC_1+kA\chi (\psi )4C_2\chi (\psi )\right)|u|^2.$$
For sufficiently $`A`$ and since $`\chi 4`$ we derive easily an estimate analogous to (2.1). From this point the proof of Theorem 2.1 applies whith just notational changes. ∎
## §3 Coverings of some strongly pseudoconcave manifolds
Let us recall the solution of the Grauert-Riemenschneider conjecture (\[GR\], p. 277) as given by Siu \[Si\] and Demailly \[De1\]. Namely the Siu–Demailly criterion says that if $`X`$ be a compact complex manifold and $`E`$ a line bundle over $`X`$. and either $`E`$ is semi-positive and positive at one point (Siu’s criterion), or
$$_{X(1)}\left(ı𝚌(E)\right)^n>0$$
$`(D)`$
(Demailly’s criterion) then $`dimH^0(X,E^k)k^n`$, for large $`k`$ and $`X`$ is Moishezon. Our aim is to extend this result in two directions. We allow $`X`$ to belong to certain classes of strongly pseudoconcave manifolds and we study (directly) Galois coverings of such manifolds.
For $`1`$– concave and compact manifolds (all which are pseudoconcave in the sense of Andreotti \[An\]) the transcendence degree of the meromorphic function field is less than or equal to the dimension of $`X`$. In the latter case we say that the manifold is Moishezon by extending the terminology from compact manifolds.
If, in the Andreotti–Grauert definition, the function $`\phi `$ can be taken such that $`a=inf\phi =\mathrm{}`$, we say that $`X`$ is hyper $`1`$– concave. Let us note that not all $`1`$– concave manifolds are hyper $`1`$– concave. Indeed, the complement of $`S^1^1`$ in $`^1`$ is $`1`$– concave but cannot possibly be hyper $`1`$– concave since $`S^1`$ is not a polar set in $``$ (I have learnt this example from M. Colţoiu and V. Vâjâitu).
Let us describe some examples. Let $`Y`$ be a compact complex manifold, $`S`$ a complete pluripolar set (the set where a strictly psh function takes the value $`\mathrm{}`$). Then $`M=YS`$ is hyper 1– concave. Conversely, if $`dimM3`$ any hyper 1– concave manifold $`M`$ is biholomorphic to a complement of a pluripolar set in a compact manifold as a consequence of Rossi’s compactification theorem. Another example of hyper $`1`$– concave manifold is $`\mathrm{Reg}(X)`$ where $`M`$ is a compact complex space with isolated singularities. Suppose that $`p`$ is an isolated singular point and that the germ $`(X,p)`$ is embedded in the germ $`(^N,0)`$ and $`z=(z_1,\mathrm{},z_N)`$ are local coordinates in the ambient space $`^N`$. The function $`\phi `$ is then obtained by cutting-off functions of the type $`\mathrm{log}(|z|^2)`$. If $`M`$ is a complete Kähler manifold of finite volume and bounded negative sectional curvature, $`M`$ is hyper $`1`$– concave. This is shown by Siu–Yau in \[SY\] by using Buseman functions. Moreover, if $`dimM3`$, this example falls in the previous case since by \[Nad\] $`M`$ can be compactified to an algebraic space by adding finitely many points.
###### Theorem 3.1
Let $`M`$ be a hyper 1– concave manifold carrying a line bundle $`(E,h)`$ which is semi-positive outside a compact set. Let $`\stackrel{~}{M}`$ be a Galois covering of group $`\mathrm{\Gamma }`$ and $`\stackrel{~}{E}`$ the lifting of $`E`$. Then
$$dim_\mathrm{\Gamma }H_{(2)}^{n,0}(\stackrel{~}{M},\stackrel{~}{E}^k)\frac{k^n}{n!}_{M(1,h)}\left(\frac{ı}{2\pi }𝚌(E,h)\right)^n+o(k^n),k\mathrm{},$$
where the $`L^2`$ condition is with respect to $`\stackrel{~}{h}`$ and any metric on $`\stackrel{~}{M}`$.
###### Demonstration Proof
Let us consider a proper function $`\phi :M(\mathrm{},0)`$ which is strictly plurisubharmonic outside a compact set. The fact that $`\phi `$ goes to $`\mathrm{}`$ to the ideal boundary of $`M`$ allows to construct a complete hermitian metric on $`M`$ which has moreover the feature of being Kähler outside a compact set. Namely we consider the function $`\chi =\mathrm{log}(\phi )`$ so that $`\overline{}\chi =\phi ^2\phi \overline{}\phi \phi ^1\overline{}\phi `$ which is obviously positive definite on the set where $`\overline{}\phi `$ is. We can now patch $`\overline{}\chi `$ and an arbitrary hermitian metric on $`M`$ by using a smooth partition of unity to get a metric $`\omega _0`$ on $`M`$ such that $`\omega _0=\overline{}\chi \text{on}MK,KM`$. It is easy to verify that $`\omega _0`$ is complete since the function $`\chi `$ is an exhaustion function and $`\omega _0=\omega +(\chi )\overline{}(\chi )`$ where $`\omega =\phi ^1\overline{}\phi `$ is a metric on $`MK`$, so that $`d(\chi )`$ is bounded in the metric $`\omega _0`$ . Note that $`\omega _0`$ is obviously Kähler on $`MK`$.
Let us consider a holomorphic hermitian line bundle $`E`$ endowed with a metric $`h`$ such that $`ı𝚌(E,h)0`$ on $`MK`$ (we stretch $`K`$ if necessary). We equip $`E`$ with the metric $`h_\epsilon =h\mathrm{exp}(\epsilon \chi )`$ and the curvature relative to the new metric satisfies $`ı𝚌(E,h_\epsilon )\epsilon \omega _0`$ on $`MK`$. We are therefore in the conditions of Theorem 2.1. First observe that $`h_\epsilon h`$ so that $`H_{(2)}^{n,0}(\stackrel{~}{M},\stackrel{~}{E}^k,\stackrel{~}{\omega }_0,\stackrel{~}{h}_\epsilon )H_{(2)}^{n,0}(\stackrel{~}{M},\stackrel{~}{E}^k,\stackrel{~}{\omega }_0,\stackrel{~}{h})`$ which is an injective $`\mathrm{\Gamma }`$–morphism. By Theorem 2.1
$$\underset{k}{lim\; inf}k^ndim_\mathrm{\Gamma }H_{(2)}^{n,0}(\stackrel{~}{M},\stackrel{~}{E}^k,\stackrel{~}{\omega }_0,\stackrel{~}{h}_\epsilon )\frac{1}{n!}_{\mathrm{\Omega }(1,h_\epsilon )}\left(\frac{ı}{2\pi }𝚌(E,h_\epsilon )\right)^n$$
so that
$$\underset{k}{lim\; inf}k^ndim_\mathrm{\Gamma }H_{(2)}^{n,0}(\stackrel{~}{M},\stackrel{~}{E}^k)\frac{1}{n!}_{\mathrm{\Omega }(1,h_\epsilon )}\left(\frac{ı}{2\pi }𝚌(E,h_\epsilon )\right)^n$$
$`(4.1)`$
We let now $`\epsilon 0`$ in (4.1); since $`h_\epsilon `$ converges uniformly together with its derivatives to $`h`$ on compact sets we see that we can replace $`h_\epsilon `$ with $`h`$ in the right-hand side of (4.1). Let $`M(q,h)`$ be the set where $`ı𝚌(E,h)`$ is non-degenerate and has exactly $`q`$ negative eigenvalues. By hypothesis $`M(1,h)K`$ and on $`M(0,h)=M(1,h)M(1,h)`$ the integrand is positive. Hence we can let $`\mathrm{\Omega }`$ exhaust $`X`$ and we get the inequality from the statement of the theorem. ∎
We prove now that Siu’s criterion extends tale quale for hyper 1–concave manifolds.
###### Corollary 3.2
Let $`M`$ be a hyper 1– concave manifold carrying a line bundle which is semi-positive outside a compact set and satisfies Demailly’s condition rm (D). Then $`X`$ is Moishezon. In particular the conclusion holds true if $`E`$ is semi-positive and positive at one point.
###### Demonstration Proof
By Theorem 3.1 (for $`\mathrm{\Gamma }=\{\mathrm{Id}\}`$) we have
$$dimH^0(M,E^kK_M)dimH_{(2)}^{n,0}(M,E^k)Ck^n$$
with $`C>0`$ for large $`k`$, by condition (D). We note that the first space is finite dimensional since $`M`$ is $`1`$–concave. By the Siegel–Serre Lemma (Proposition 5.7 from \[Ma\]), $`dimH^0(M,E^kK_M)Ck^{\varkappa (E)}`$, ($`k>0`$), where $`\varkappa (E)`$ is the supremum over $`k`$ of the generic rank of the canonical meromorphic mapping from $`M`$ to $`\left(H^0(M,E^kK_M)^{}\right)`$. We obtain that $`\varkappa (E)=n`$, that is, the line bundle $`E^kK_M`$ gives local coordinates on an open dense set of $`M`$ for sufficiently large $`k`$. This clearly implies $`M`$ Moishezon and thereby concludes the proof. ∎
###### Remark Remark 3.1
(a) We can use this criterion in the Nadel compactification theorem \[Nad\]. It asserts that if $`M`$ is a connected manifold of dimension $`3`$ satisfying : (i) $`M`$ is hyper 1–concave, (ii) $`M`$ is Moishezon, (iii) $`M`$ can be covered by Zariski-open sets which are uniformized by Stein manifolds, then $`M`$ is biholomorphic to $`M^{}S`$ where $`M^{}`$ is a compact Moishezon space and $`S`$ is finite. We see thus that condition (ii) in Nadel’s theorem may be replaced with the analytic condition: $`M`$ possesses a line bundle which is semi-positive outside a compact set and satisfies Demailly’s condition (D).
(b) In general, if $`M`$ is a hyper $`1`$– convave manifold of dimension $`n3`$ possesing a semi–positive line bundle satistying (D) then (by a theorem of Rossi) it can be compactified so that $`M`$ is biholomorphic to an open set of a compact Moishezon manifold which is the complement of a complete pluripolar set. Therefore there exist a meromorphic mapping defined on $`X`$ with values in a projective space which is an embedding outside a proper analytic set of $`X`$. To see this we have to apply the corresponding statement for compact Moishezon manifolds, a result due to Moishezon. The difficulty in Nadel’s theorem is to show that under additional hypothesis the pluripolar set is actually a finite set.
(c) The argument in the proof of Corollary 3.2 shows that the integral appearing in Theorem 3.1 is finite. Thus, if $`E`$ is positive outside a compact set $`K`$ then $`MK`$ has finite volume with respect to the metric $`ı𝚌(E)`$ (this observation stems from \[NT\]).
(d) If $`M`$ possesses a positive line bundle $`E`$ then $`ı𝚌(E)+ı\overline{}\chi `$ is a complete Kähler metric and Hörmander’s $`L^2`$ estimates and Andreotti–Tomassini’s theorem \[AT\] show that $`E`$ is ample and $`M`$ can be embedded in the projective space. So even in dimension $`2`$ we can compactify $`M`$ (by \[An\]).
(f) Let $`X`$ is a compact complex space of dimension $`n2`$ and with isolated singularities. Suppose that we have a line bundle $`E`$ on $`\mathrm{Reg}(X)`$ which is semi-positive in a deleted neighbourhood of $`\mathrm{Sing}(X)`$ and satisfies (D). Then $`X`$ is Moishezon. Indeed, by the previous result we find $`n=dimX`$ independent meromorphic functions on $`\mathrm{Reg}(X)`$ which extend to $`X`$ by the Levi extension theorem. This is a generalization of Takayama’s criterion \[Ta\] in the case of isolated singularities. We allow weaker hypothesis, that is $`E`$ is defined just on $`\mathrm{Reg}(X)`$ and the curvature condition is just semi-positivity. The reason is the good exhaustion function we have at hand. In the general case one has to use the Poincaré metric and the strict positivity near $`\mathrm{Sing}(X)`$ is essential. Note however that the method of Takayama gives that the line bundle who forms local coordinates is $`E^k`$, while in our proof is $`E^kK_X`$.
We want now to study the following type of stongly pseudoconcave manifold. Let $`X`$ be an irreducible compact complex space with isolated singularities and of dimension $`2`$. We know that $`\mathrm{Reg}(X)`$ is hyper $`1`$– concave and we denote by $`\phi :\mathrm{Reg}(X)`$ the exhaustion function. Since $`\phi `$ is strictly plurisubharmonic outside a compact set we have that the sub–level sets $`X_c=\{\phi >c\}`$ are $`1`$– concave manifolds i.e. stongly pseudoconcave domains. In our previous paper \[Ma\] we have shown that in general if $`M`$ is a $`1`$–concave manifold of dimension $`3`$ which carries a hermitian line bundle $`E`$ which semi-negative near the boundary and satisfies (D) then the Kodaira dimension of $`E`$ is maximal and $`M`$ is Moishezon. The assumption about the change of curvature sign (i.e. semi-negativity) near the boundary is imposed by the construction of complete hermitian metrics $`\omega _0`$ and $`h_0`$ as in Theorem 2.2 which give the $`L^2`$ estimate and preserve condition (D) for $`h_0`$; the negativity of the Levi form of the sublevel sets of $`M`$ requires as a natural curvature condition for $`E`$ the semi–negativity. The restriction on dimension comes from the fact that we need an $`L^2`$ estimate in bi–degree $`(0,1)`$. Of course, usually we are given an overall positive bundle $`E`$ on $`M`$. We show that for manifolds $`X_c`$ as before we can also apply the criterion in \[Ma\] alluded to by modifying the metric.
We recall at the outset some terminology. Let us consider a covering $`\{U_\alpha \}`$ of $`X`$ and embeddings $`\iota _\alpha :U_\alpha ^{N_\alpha }`$ such that $`E|_{U_\alpha }`$ is the inverse image by $`\iota _\alpha `$ of the trivial line bundle $`_\alpha `$ on $`^{N_\alpha }`$. Moreover we consider hermitian metrics $`h_\alpha =e^{\phi _\alpha }`$ on $`_\alpha `$ such that $`\iota _\alpha ^{}h_\alpha =\iota _\beta ^{}h_\beta `$ on $`U_\alpha U_\beta \mathrm{Reg}(X)`$. The system $`h=\{\iota _\alpha ^{}h_\alpha \}`$ is called a hermitian metric on $`E`$ over $`X`$. It clearly induces a hermitian metric on $`E`$ over $`\mathrm{Reg}(X)`$. The curvature current $`ı𝚌(E)`$ is given in $`U_\alpha `$ by $`\iota _\alpha ^{}(ı\overline{}\phi _\alpha )`$ which on $`\mathrm{Reg}(X)`$ agrees with the curvature of the induced metric.
###### Theorem 3.3
Let $`X`$ be an irreducible compact complex space with isolated singularities and let $`X_c`$ be the sublevel sets of the hyper $`1`$– concave manifold $`\mathrm{Reg}(X)`$. Assume that there exists a holomorphic line bundle $`EX`$ with a smooth hermitian metric such that condition rm (D) is fulfilled on $`\mathrm{Reg}(X)`$. Then for sufficiently small $`c`$ there exists a metric on $`E`$ such that $`E`$ is negative in the neighbourhood of $`bX_c`$ and $`_{X_c(1)}\left(ı𝚌(E)\right)^n>0`$.
###### Demonstration Proof
Let $`\pi :\stackrel{~}{X}X`$ be a resolution of singularities of $`X`$. Let us denote by $`D_i`$ the components of the exceptional divisor. Then there exist positive integers $`n_i`$ such that $`D:=n_iD_i`$ admits a smooth hermitian metric such that the induced line bundle $`[D]`$ is negative in a neighbourhood $`\stackrel{~}{U}`$ of $`D`$ (cf. \[Sa\]). Let us consider a canonical section $`s`$ of $`[D]`$, i.e. $`D=(s)`$, and denote by $`|s|^2`$ the pointwise norm of $`s`$ with respect to the above metric. By Lelong-Poincaré equation $`\phi =\mathrm{log}|s|^2`$ is strictly plurisubharmonic on $`\stackrel{~}{U}D`$. By using a smooth function on $`\stackrel{~}{X}`$ with compact support in $`\stackrel{~}{U}`$ which equals one near $`D`$ we construct a smooth function $`\chi `$ on $`\stackrel{~}{X}D\mathrm{Reg}(X)`$ such that $`\chi =\mathrm{log}(\mathrm{log}|s|^2)`$ on $`\stackrel{~}{U}D`$.
Since $`\mathrm{log}|s|^2`$ goes to $`\mathrm{}`$ on $`D`$, this is the analogue of the function constructed in the proof of Theorem 3.1 . As there we show that $`ı\chi \overline{}\chi ı\overline{}\chi `$. Let us consider a metric $`\omega `$ on $`\mathrm{Reg}(X)`$ which on every open set $`U_\alpha `$ as above is the pullback of a hermitian metric on the ambient space $`^{N_\alpha }`$, $`\omega =\iota _\alpha ^{}\omega _\alpha `$ . We consider then the metric (Kähler near $`\mathrm{Sing}(X)`$) $`\omega _0=A\omega +\overline{}\chi `$ where $`A>0`$ is chosen sufficiently large (to ensure that $`\omega _0`$ is a metric away from the open set where $`\overline{}\chi `$ is positive definite). It is easily seen that $`\omega _0`$ is complete by the same argument as in the proof of Theorem 3.1 . This kind of metrics were introduced by Saper in \[Sa\]. They have finite volume.
Let us consider now a neighbourhood $`U`$ of the singular set. We assume that $`U`$ is small enough so that there are well defined on $`U`$ a potential $`\rho `$ for $`\omega `$ and a potential $`\eta `$ for the curvature $`ı𝚌(E)`$ (they are restrictions from ambient spaces). By suitably cutting-off we may define a function $`\psi 𝒞^{\mathrm{}}(\mathrm{Reg}(X))`$ such that $`\psi =\chi \eta A\rho `$ near $`\mathrm{Sing}(X)`$ . Remark that since $`ı𝚌(E)`$ is bounded above by a continuous $`(1,1)`$ form near $`\mathrm{Sing}(X)`$ the potential $`\eta `$ is bounded above near the singular set. This holds true for $`\rho `$ too (it is smooth) so that $`\psi `$ tends to $`\mathrm{}`$ at the singular set $`\mathrm{Sing}(X)`$. Let us consider a smooth function $`\gamma :`$ such that
$$\gamma (t)=\{\begin{array}{cc}0\hfill & \text{if}t0,\hfill \\ t\hfill & \text{if}t1.\hfill \end{array}$$
and the functions $`\gamma _\nu :`$ given by $`\gamma _\nu (t)=\gamma (t\nu )`$ for all positive integers $`\nu `$ . Let us denote the hermitian metric on $`E`$ by $`h`$ and let us consider the following metric on $`E`$ : $`h_\nu =h\mathrm{exp}\left(\gamma _\nu (\psi )\right)`$, with curvature
$$ı𝚌(E,h_\nu )=ı𝚌(E,h)+\gamma _\nu ^{}(\psi )\overline{}\psi +\gamma _\nu ^{\prime \prime }(\psi )\psi \overline{}\psi .$$
On the set $`\{\psi \nu +1\}`$ we have $`\gamma _\nu (\psi )=\psi \nu `$ so that $`\gamma _\nu ^{}(\psi )=1`$ and $`\gamma _\nu ^{\prime \prime }(\psi )=0`$ and therefore $`ı𝚌(E,h_\nu )=ı𝚌(E,h)+\overline{}\psi `$. Since $`\psi `$ goes to $`\mathrm{}`$ when we approach the singular set we may choose $`\nu _0`$ such that for $`\nu \nu _0`$ we have $`\{\psi \nu +1\}U`$ where $`U`$ is a sufficiently small neighbourhood of $`\mathrm{Sing}(X)`$. Bearing in mind the meaning of $`\eta `$ and $`\rho `$ together with the definition of $`\omega _0`$ it is straightforward that $`ı𝚌(E,h_\nu )=\omega _0`$ on $`\{\psi \nu +1\}`$, that is $`(E,h_\nu )`$ is negative on this set. We denote $`\mathrm{\Omega }_\nu `$ the compact set $`\{\psi \nu +2\}`$ . We decompose this set in $`\mathrm{\Omega }_\nu ^{}=\{\psi \nu \}`$ and $`\mathrm{\Omega }_\nu ^{\prime \prime }=\{\nu \psi \nu +2\}`$. On $`\mathrm{\Omega }_\nu ^{}`$ we have $`\gamma _\nu (\psi )=0`$ and $`ı𝚌(E,h_\nu )=ı𝚌(E,h)`$ . We infer that
$$\begin{array}{ccc}\hfill _{\mathrm{\Omega }_\nu ^{}(1,h_\nu )}\left(ı𝚌(E,h_\nu )\right)^n& =_{\mathrm{\Omega }_\nu ^{}(1,h)}\left(ı𝚌(E,h)\right)^n\hfill & \\ & =_{\mathrm{Reg}(X)(1,h)}\mathrm{𝟷}_{\mathrm{\Omega }_\nu ^{}}\alpha _1\mathrm{}\alpha _n𝑑V_0\hfill & (4.2)\hfill \end{array}$$
where $`\alpha _1,\mathrm{},\alpha _n`$ are the eigenvalues of $`ı𝚌(E,h)`$ with respect to $`\omega _0`$ and $`dV_0`$ is the volume form of the same metric. Since $`ı𝚌(E,h)`$ is dominated by the euclidian metric near $`\mathrm{Sing}(X)`$, $`ı𝚌(E,h)`$ is dominated by $`\omega `$ and by $`\omega _0`$. Hence the product $`\alpha _1\mathrm{}\alpha _n`$ is bounded on $`\mathrm{Reg}(X)`$. Since $`\mathrm{Reg}(X)(1)`$ has finite volume with respect to $`\omega _0`$ the functions $`|\mathrm{𝟷}_{\mathrm{\Omega }_\nu ^{}}\alpha _1\mathrm{}\alpha _n|`$ are bounded by an integrable function. On the other hand $`\mathrm{𝟷}_{\mathrm{\Omega }_\nu ^{}}1`$ when $`\nu \mathrm{}`$ so that the integrals in (4.2) tend to $`_{\mathrm{Reg}(X)(1,h)}\left(ı𝚌(E,h)\right)^n`$ which is assumed to be positive. Thus it suffices to show that the integral on the set $`\mathrm{\Omega }_\nu ^{\prime \prime }`$ i.e. $`_{\mathrm{\Omega }_\nu ^{\prime \prime }(1,h_\nu )}\left(ı𝚌(E,h_\nu )\right)^n`$ tends to zero as $`\nu \mathrm{}`$. For this purpose we use the obvious bound
$$_{\mathrm{\Omega }_\nu ^{\prime \prime }(1,h_\nu )}\left(\frac{ı}{2\pi }𝚌(E,h_\nu )\right)^nsup|\delta _1\mathrm{}\delta _n|\mathrm{vol}(\mathrm{\Omega }_\nu ^{\prime \prime })$$
where $`\delta _1,\mathrm{},\delta _n`$ are the eigenvalues of $`ı𝚌(E,h_\nu )`$ with respect to $`\omega _0`$ and the volume is taken in the same metric. We use now the minimum-maximum principle to see that: (i) $`\delta _1`$ is bounded below and $`\delta _2,\mathrm{},\delta _n`$ are bounded above on the set of integration $`\mathrm{\Omega }_\nu ^{\prime \prime }(1,h_\nu )`$ and (ii) $`\delta _1,\mathrm{},\delta _n`$ are upper bounded on $`\mathrm{\Omega }_\nu ^{\prime \prime }(0,h_\nu )`$. For this we need the domination of $`ı𝚌(E,h)`$ by $`\omega `$ and the boundedness of $`\gamma _\nu ^{}`$ and $`\gamma _\nu ^{\prime \prime }`$ . Since $`\mathrm{vol}(\mathrm{\Omega }_\nu ^{\prime \prime })0`$ as $`\nu \mathrm{}`$ our contention follows. Hence for large $`\nu `$ the metric $`h_\nu `$ does the required job. ∎
###### Remark Remark 3.2
We have seen that Siu’s criterion generalizes to compact complex spaces with isolated singularities. Demailly’s criterion extends too. Let $`X`$ be an irreducible compact complex space of dimension $`n3`$ with isolated singularities and $`E`$ a smooth hermitian line bundle over $`X`$. Assume that condition rm (D) is fulfilled on $`\mathrm{Reg}(X)`$. Then $`X`$ is Moishezon. Indeed, for small $`c`$ the sets $`X_c`$ are Moishezon by Corollary 4.3 of \[Ma\] and the meromorphic functions from $`X_c`$ extend to $`X`$. In fact the result holds also for $`n=2`$ with a proof very similar to that of Theorem 4.4. We note also that we can allow the metric $`h`$ of $`E`$ to be singular at $`\mathrm{Sing}(X)`$ but the cuvature current $`ı𝚌(E)`$ should be dominated (above ans below) by the euclidian metric near $`\mathrm{Sing}(X)`$. The proof of Theorem 4.4 goes through with minor changes.
Since the manifold $`\overline{X}_c`$ is compact Theorem 4.4 can be used to prove some stability results for certain perturbation of the complex structure of $`\overline{X}_c`$. Since our approach relies on the use of a sufficiently positive line bundle $`E`$ we need to consider perturbations of the complex structure which lift to a perturbation of $`E`$. This kind of sitiuation was studied by L. Lempert in \[Le\].
###### Proposition 3.4
Let $`X`$ be a Moishezon variety with isolated singularities and dimension $`n3`$. Let $`𝒥`$ denote the complex structure of $`\mathrm{Reg}(X)`$ and let $`Z\mathrm{Reg}(X)`$ be a non–singular hypersurface such that the line bundle $`E=[Z]`$ satisfies rm (D). Then for sufficiently small $`c`$ and any complex structure $`𝒥^{}`$ on $`\overline{X}_c`$ such that
there exists a $`𝒥^{}`$–holomorphic line bundle $`E^{}`$ on $`\overline{X}_c`$ which is negative near $`bX_c`$ and satisfies rm (D). In particular $`(\overline{X}_c,𝒥^{})`$ is a Moishezon pseudoconcave manifold and any compactification of $`(\overline{X}_c,𝒥^{})`$ is Moishezon.
###### Demonstration Proof
Let us first choose $`c_0`$ such that for $`c<c_0`$ there exists a ‘good’ hermitian metric $`h`$ on $`E`$ over a neighbourhood of $`X_c`$, that is, whith negative curvature near the boundary and satisfying (D). We use now the description of the lifting of $`𝒥^{}`$ with properties (1) and (2) as given in \[Le\]. Namely, $`Z`$ determines a new $`𝒥^{}`$ holomorphic line bundle $`E^{}(\overline{X}_c,𝒥^{})`$. There exists a finite open covering $`𝒰=\{U\}`$ of $`\overline{X}_c`$ such that $`E`$ and $`E^{}`$ are trivial on each $`U`$ and they are defined by multiplicative cocycles $`\{g_{UV}𝒥\text{holomorphic on}\overline{U}\overline{V}:U,V𝒰\}`$ and $`\{g_{UV}^{}𝒥^{}\text{holomorphic on}\overline{U}\overline{V}:U,V𝒰\}`$. Moreover $`g_{UV}`$ and $`g_{UV}^{}`$ are as close as we please assuming $`𝒥`$ and $`𝒥^{}`$ are sufficiently close. (By ‘close’ we always understand close in the $`𝒞^{\mathrm{}}`$ topology.) Next we can define a smooth bundle isomorphism $`EE^{}`$ by resolving the smooth additive cocycle $`\mathrm{log}(g_{UV}^{}/g_{UV})`$ in order to find smooth functions $`f_U`$, close to 1 on a neighbourhood of $`\overline{U}`$ such that $`g_{UV}^{}=f_Ug_{UV}f_V^1`$. Then the isomorphism between $`E`$ and $`E^{}`$ is given by $`f=\{f_U\}`$. The metric $`h`$ is given in terms of the covering $`𝒰`$ by a collection $`h=\{h_U\}`$ of smooth strictly positive functions satisfying the relation $`h_V=h_U|g_{UV}|`$. We define a hermitian metric $`h^{}=\{h_U^{}\}`$ on $`E^{}`$ by $`h_U^{}=h_U|f_U^1|`$; $`h_U^{}`$ is close to $`h_U`$. The curvatures forms of $`E`$ and $`E^{}`$ are given by
$$\frac{ı}{2\pi }𝕔(E)=\frac{1}{4\pi }d𝒥d(\mathrm{log}h_U),\frac{ı}{2\pi }𝕔(E^{})=\frac{1}{4\pi }d𝒥^{}d(\mathrm{log}h_U^{}).$$
Therefore, when $`𝒥^{}`$ is sufficiently close to $`𝒥`$, $`\frac{ı}{2\pi }𝕔(E^{})`$ is negative near the boundary of $`\overline{X}_c`$ and, since the eigenvalues of $`\frac{ı}{2\pi }𝕔(E^{})`$ are close to those of $`\frac{ı}{2\pi }𝕔(E)`$, $`E^{}`$ satisfy the condition (D) i.e. $`_{X_c(1)}\left(ı𝚌(E^{})\right)^n>0`$. We can apply thus the Corrolary 4.3 of \[Ma\] to the strongly pseudoconcave manifold $`(\overline{X}_c,𝒥^{})`$ to conclude that $`(\overline{X}_c,𝒥^{})`$ is Moishezon. ∎
###### Remark Remark 3.3
If $`[Z]`$ is positive, part of the stability property follows from the rigidity of embeddings with positive normal bundle. Indeed, assume $`N_Z=[Z]_Z`$ is positive in $`(\overline{X}_c,𝒥^{})`$ (for any $`c`$ such that this manifold is still pseudoconcave). Then Ph. Griffiths \[Gri1\] has shown that there exists a neighbourhood $`W`$ of $`Z`$ such that the mapping $`\mathrm{\Phi }:(X_c,𝒥^{})^N`$ given by $`[mZ]`$ is an embedding of $`W`$ for large $`m`$ . Thus $`(X_c,𝒥^{})`$ is Moishezon. Our result deals with the slightly more general situation of a ‘big’ embedding i.e. when $`[Z]`$ is not ample but satisfies condition (D). Moreover we have a useful quantitative way of measuring whether the perturbed structure is Moishezon.
###### Corollary 3.5
Let $`(\overline{X}_c,𝒥^{})`$ and $`E^{}`$ be as in Proposition 4.6. Then there exists hermitian metrics on $`X_c`$ and $`E^{}`$ and a positive constant $`C`$ such that for any Galois covering $`\stackrel{~}{X}_cX_c`$ of group $`\mathrm{\Gamma }`$ we have
$$dim_\mathrm{\Gamma }H_{(2)}^0(\stackrel{~}{X}_c,\stackrel{~}{E^{}}^k)Ck^n+o(k^n),k\mathrm{}.$$
the $`L^2`$ condition being with respect to lifts of the hermitian metrics on $`X_c`$ and $`E^{}`$.
###### Demonstration Proof
We know that we have on $`E^{}`$ a metric $`h`$ satisfying the conclusion of Theorem 4.4. Then, as in Theorem 2.2, we can construct metrics $`\omega _0`$ and $`h_0`$ in order to obtain (2.9). Note that the integral in (2.9) depends on the modified metric $`h_0`$ so we cannot always infer that it is positive even if $`(E^{},h)`$ satisfies (D). But under the assumption of semi–negativity of $`h`$ near the boundary we can construct an $`h_0`$ such that the integral in (2.9) is positive (cf. Corollary 4.3 of \[Ma\]). Thus by applying Theorem 2.2 we get the conclusion. ∎
## §4 $`L^2`$ generalization of a theorem of Takayama
In this section we study the $`L^2`$ cohomology of coverings of Zariski open sets in compact complex spaces. For compact spaces with singularities Takayama \[Ta\] generalized Siu–Demailly criterion if $`EX`$ is a line bundle endowed with a singular hermitian metric which is smooth outside a proper analytic set $`Z\mathrm{Sing}(X)`$ and defines a strictly positive current near $`Z`$.
Using the setting of Takayama’s theorem we shall study coverings of Zariski open sets in compact complex spaces.
###### Proposition 4.1
Let $`X`$ be an $`n`$–dimensional compact manifold and let $`E`$ be a holomorphic line bundle with a singular hermitian metric $`h`$. We assume that:
Let $`p:\stackrel{~}{M}M`$ be a Galois covering with group $`\mathrm{\Gamma }`$ and $`\stackrel{~}{E}=p^{}E`$. Then,
$$dim_\mathrm{\Gamma }H_{(2)}^0(\stackrel{~}{M},\stackrel{~}{E}^k)\frac{k^n}{n!}_{M(1,h)}\left(\frac{ı}{2\pi }𝚌(E,h)\right)^n+o(k^n),k>>0,$$
where $`H_{(2)}^0(\stackrel{~}{M},\stackrel{~}{E}^k)`$ is the space of sections of $`\stackrel{~}{E}^k`$ which are $`L^2`$ with respect to the pullbacks of the restrictions to $`M`$ and $`E_M`$ of smooth metrics on $`X`$ and $`E`$.
###### Demonstration Proof
This is an equivariant form of Takayama’s main technical result in \[Ta\]. Namely we construct the Poincaré metric $`\omega _\epsilon `$ on $`M`$ (for details see \[Zu\]) and $`h_\epsilon `$ as in \[Ta\] and remark that the hypothesis of Theorem 2.1 are satisfied. Moreover we can work with $`(0,1)`$–forms since the Ricci curvature of the Poincaré metric is bounded below.
More specifically, we write $`Z=Z_j`$ and consider a section $`\sigma _j`$ of the line bundle $`[Z_j]`$ which vanishes to first order on $`Z_j`$. Then we endow $`[Z_j]`$ with a hermitian metric such that the norm of $`\sigma _j`$ satisfies $`|\sigma _j|<1`$. Take then an arbitrary smooth metric $`\omega ^{}`$ on $`X`$ and define $`\omega _\epsilon =\omega ^{}\epsilon ı\overline{}(\mathrm{log}|\sigma _j|^2)^2`$ on $`M=XZ`$ which for small $`\epsilon >0`$ is a complete metric on $`M`$. Then we consider the following family of metrics on $`E_M`$: $`h_\epsilon =h_j(\mathrm{log}|\sigma _j|^2)^\epsilon `$, $`\epsilon >0`$. We check now the hypotheses of Theorem 2.1 is satisfied. First we remark that the torsion operators of the Poincaré metric are pointwise bounded with respect to the Poincaré metric since $`d\omega _\epsilon =d\omega ^{}`$ and $`\omega _\epsilon >\omega ^{}`$. Also the Ricci curvature $`𝚌(K_{\stackrel{~}{M}}^{})`$ of $`\stackrel{~}{\omega }_\epsilon `$ is bounded below with respect to $`\stackrel{~}{\omega }_\epsilon `$ by a constant independent of $`\epsilon `$ (since this is true for $`\omega _\epsilon `$). Since $`E`$ is strictly positive in the neighbourhood of $`Z`$ condition (A) is satisfied for a compact $`K`$ outside which $`E`$ is positive (and it doesn’t depend on $`\epsilon `$).
Let $`h^{}`$ be a smooth hermitian metric on $`E`$ over $`X`$. Near $`Z`$ the metric $`h`$ is locally represented by a strictly plurisubharmonic weight. Thus $`h`$ is locally bounded below near $`Z`$ and thus $`hCh^{}`$ on $`X`$ for some positive constant $`C`$. We remark now that $`h_\epsilon >hCh^{}`$ and $`\omega _\epsilon >\omega ^{}`$ near $`Z`$ so that we have the inclusion $`H_{(2)}^0(\stackrel{~}{M},\stackrel{~}{E}^k)_\epsilon H_{(2)}^0(\stackrel{~}{M},\stackrel{~}{E}^k)`$, (which is an injective $`\mathrm{\Gamma }`$–morphism) in the last group the $`L^2`$ condition being taken with respect to $`\stackrel{~}{h}^{}`$ and $`\stackrel{~}{\omega }^{}`$. By Theorem 2.1 for $`K\mathrm{\Omega }M`$
$$dim_\mathrm{\Gamma }H_{(2)}^0(\stackrel{~}{M},\stackrel{~}{E}^k)dim_\mathrm{\Gamma }H_{(2)}^0(\stackrel{~}{M},\stackrel{~}{E}^k)_\epsilon _{\mathrm{\Omega }(1,h_\epsilon )}\left(\frac{ı}{2\pi }𝚌(E,h_\epsilon )\right)^n+o(k^n).$$
We can let $`\epsilon 0`$ in the right–hand side in order to replace $`h_\epsilon `$ with $`h`$. Then we can let $`\mathrm{\Omega }`$ exhaust $`X`$ to get the inequality from the statement. ∎
###### Theorem 4.2
Let $`X`$ be an irreducible reduced compact Moishezon space and let $`M\mathrm{Reg}(X)`$ be a Zariski open set. There exists a holomorphic line bundle $`E\mathrm{Reg}(X)`$ endowed with a singular hermitian metric whose curvature current $`ı𝚌(E)`$ is positive and such that for any Galois covering $`\stackrel{~}{M}\mathrm{@}>p>>M`$ of group $`\mathrm{\Gamma }`$ we have
$$dim_\mathrm{\Gamma }H_{(2)}^0(\stackrel{~}{M},\stackrel{~}{E}^k)\frac{k^n}{n!}_M\left(\frac{ı}{2\pi }𝚌(E)\right)^n+o(k^n),k\mathrm{}$$
where the integration takes place outside $`\mathrm{Sing}\mathrm{supp}𝚌(E)`$. The $`L^2`$ condition is taken with respect to liftings of smooth hermitian metrics on $`M`$ and $`E`$ induced from a resolution of singularities of $`X`$.
###### Demonstration Proof
Step 1. Let $`X`$ be a Moishezon manifold and $`M=XZ`$ a Zariski open set, where $`Z`$ is a proper analytic set. Thanks to Moishezon $`X`$ admits a projective modification. Therefore there exists a strictly positive integral Kähler current $`T`$ on $`X`$. Equivalently there exists a holomorphic line bundle $`E`$ on $`X`$ possesing a singular hermitian metric such that the curvature current $`T=ı𝚌(E)`$ is strictly positive (bounded below by a smooth hermitian metric). Assume that $`\mathrm{Sing}\mathrm{supp}TZ`$. Then $`M`$ is biholomorphic to a Zariski open set as in the statement of Proposition 3.1. Indeed, we can blow up $`Z`$ to make it a divisor with only simple normal crossings. By replacing $`E`$ with higher tensor powers and twisting it with the dual of the exceptional divisor at each step of the blowing up process we can ensure that on the blow–up we still have a positive line bundle with singular metric along $`Z`$. Thus in this case we can apply Proposition 3.1.
Step 2. To go further let $`M`$ be a Zariski open set in a Moishezon manifold $`X`$. By a theorem of Demailly \[De2\] we know that there exists a strictly positive integral Kähler current $`T`$ with analytic singularities. As a consequence $`\mathrm{Sing}\mathrm{supp}TS`$, where $`S`$ is a proper analytic set. As before we can suppose that $`SZ`$ is a divisor with only simple normal crossings. Let $`E`$ be a line bundle with singular hermitian metric such that $`T=ı𝚌(E)`$. Denote by $`M_1=X(SZ)=MS`$ : $`M_1`$ and $`E`$ are as in Proposition 3.1. Let $`p:\stackrel{~}{M}M`$ be a Galois covering of group $`\mathrm{\Gamma }`$. Setting $`\stackrel{~}{M}_1=p^1M_1`$ we have a Galois covering $`\stackrel{~}{M}_1M_1`$ of group $`\mathrm{\Gamma }`$. Hence, $`dim_\mathrm{\Gamma }H_{(2)}^0(\stackrel{~}{M}_1,\stackrel{~}{E}^k)\frac{k^n}{n!}_{M_1}\left(\frac{ı}{2\pi }𝚌(E)\right)^n+o(k^n)`$, for $`k\mathrm{}`$. The $`L^2`$ condition on $`\stackrel{~}{M}_1`$ is with respect to liftings of smooth hermitian metrics on $`X`$ and $`E`$. But a holomorphic section defined outside the analytic set $`\stackrel{~}{S}=p^1S`$ which is square integrable with respect to a smooth metric on $`\stackrel{~}{M}`$ extends past $`\stackrel{~}{S}`$ as a holomorphic section on $`\stackrel{~}{M}`$. We infer $`dim_\mathrm{\Gamma }H_{(2)}^0(\stackrel{~}{M},\stackrel{~}{E}^k)\frac{k^n}{n!}_M\left(\frac{ı}{2\pi }𝚌(E)\right)^n+o(k^n)`$ the integral being taken on the smooth locus of $`ı𝚌(E)`$ The $`L^2`$ condition is taken with respect to pullbacks of smooth metrics on $`X`$ and $`E`$.
Step 3. Finally let $`X`$ and $`M`$ as in hypothesis. By a resolution of singularities $`M`$ is biholomorphic to a Zariski open set of a Moishezon manifold. By the preceding remarks we can conclude. ∎
The following Proposition is a consequence of Theorem 4.2 in the case of Galois coverings (taking into acount that the number of sheets of such a covering equals the cardinal of $`\mathrm{\Gamma }`$). However, using Theorem 2.2 of Napier and Ramachandran \[NR\], we can prove it for any unramified covering.
###### Proposition 4.3
Let $`X`$ be an irreducible reduced compact Moishezon space and let $`M\mathrm{Reg}(X)`$ be a Zariski open set. There exists a holomorphic line bundle $`E\mathrm{Reg}(X)`$ such that for any unramified covering $`p:\stackrel{~}{M}M`$ we have
$$dimH_{(2)}^0(\stackrel{~}{M},\stackrel{~}{E}^k)Ck^nd,k>>0$$
$`(4.1)`$
where $`d`$ is the number of sheets of the covering and $`C>0`$.
###### Demonstration Proof
In the situation of Step 1 of the preceding proof we see that the Poincaré metric on $`M`$ is a complete Kähler metric since $`M`$ has the Kähler metric $`ı𝕔(E)`$. Therefore $`\stackrel{~}{M}`$ possesses a complete Kähler metric and a positive line bundle $`\stackrel{~}{E}`$. By applying the $`L^2`$ estimates of Hörmander as in \[NR, Theorem 2.2\] we get the result for the $`L^2`$ cohomology with respect to the metrics $`\stackrel{~}{\omega }_\epsilon `$ and $`\stackrel{~}{h}_\epsilon `$ (notations of Proposition 4.1). As in the proof of Proposition 4.1 we see that we can actually use pull-backs of smooth metrics on $`X`$. Steps 2 and 3 go through as before. ∎
## §5 Further remarks
We will apply Theorem 2.1 to the case of a complete Kähler manifold $`M`$ with positive canonical bundle $`K_M`$. The case $`\mathrm{\Gamma }=\{\mathrm{Id}\}`$ is due to Nadel and Tsuji \[NT\]. If $`D`$ is a bounded domain of holomorphy in $`^n`$ we know by a theorem of Bremermann that the Bergman metric $`\omega =\omega _B`$ is complete. On the other hand the Bergman metric is invariant under analytic automorphisms. Thus this metric descends to a complete Kähler metric on any quotient of the domain by a properly discontinuous discrete group $`\mathrm{\Gamma }\mathrm{Aut}(D)`$. We denote $`M=D/\mathrm{\Gamma }`$ and $`\omega _{}`$ the induced Bergman metric on $`M=D/\mathrm{\Gamma }`$. If we denote by $`B(z,\overline{z})`$ the Bergman kernel of $`D`$ we know that $`B^1`$ can be considered as a hermitian $`\mathrm{\Gamma }`$–invariant metric on $`K_D`$. Since $`\omega =\overline{}\mathrm{log}B(z,\overline{z})`$ there exists a hermitian metric on $`K_M`$ such that $`𝚌(K_M)=\omega _{}`$. We have thus the following.
###### Proposition 5.1
Let $`D`$ is a bounded domain of holomorphy in $`^n`$, $`\mathrm{\Gamma }\mathrm{Aut}D`$ a discrete group acting properly discontinuously on $`D`$ and $`M=D/\mathrm{\Gamma }`$. Then
$$dim_\mathrm{\Gamma }H_2^0(D,K_D^k)\left(\frac{k}{2\pi }\right)^n_M\frac{\omega _{}^n}{n!}+o(k^n),k\mathrm{}$$
where the $`L^2`$ condition is taken with respect to the Bergman metric on $`D`$ and the metric $`B^1`$ on $`K_D`$.
Note that the space $`H_2^0(D,K_D^k)`$ is a space of square integrable functions with respect to the Bergman metric and to the weight $`B^k`$. An immediate consequence is the following.
###### Corollary 5.2
Assume that the Bergman metric on $`M`$ has infinite volume. Then $`dim_\mathrm{\Gamma }H_2^0(D,K_D^k)=\mathrm{}`$ for $`k`$ large enough.
We remark that the last conclusion is stronger than the results coming from the $`L^2`$ method which gives just $`dimH_2^0(D,K_D^k)C|\mathrm{\Gamma }|k^n`$ for some positive constant $`C`$.
Let us see what become our results in the simplest case of the unit disk $`D`$. Then the Bergman metric equals the hyperbolic metric $`(1|z|^2)^2dzd\overline{z}`$. If $`\mathrm{\Gamma }`$ is a Fuchsian group, we have the following possibilities for large $`k`$:
(a) If $`M=D/\mathrm{\Gamma }`$ is compact, $`dim_\mathrm{\Gamma }H_2^0(D,K_D^k)=k\mathrm{vol}(M)+o(k)`$.
(b) If $`M`$ is non–compact and has a finite number of cusps, the hyperbolic volume $`\mathrm{vol}(M)`$ is finite and $`dim_\mathrm{\Gamma }H_2^0(D,K_D^k)k\mathrm{vol}(M)+o(k)`$,
(c) If $`M`$ is non–compact and the discontinuity set $`\mathrm{\Omega }S^1`$ is a union of intervals, $`dim_\mathrm{\Gamma }H_2^0(D,K_D^k)=\mathrm{}`$ (since $`\mathrm{vol}(M)=\mathrm{}`$).
According to a conjecture of Griffiths \[Gri2, p.50\], if $`D`$ is a bounded domain in $`^n`$ which is topologically a cell and $`D/\mathrm{\Gamma }`$ is quasi–projective then (i) the Bergman metric on $`D/\mathrm{\Gamma }`$ is complete and (ii) the volume of $`D/\mathrm{\Gamma }`$ with respect to this metric is finite. In the sequel we discuss the conjecture without the topological restriction. If $`D`$ is a domain of holomorphy and $`M=D/\mathrm{\Gamma }`$ is pseudoconcave (e.g. $`\mathrm{codim}(\overline{M}M)2`$), the answer is yes. Indeed, this follows from the Riemann–Roch inequalities for $`\mathrm{\Gamma }=\{\mathrm{Id}\}`$ in Proposition 5.1. If $`D`$ is not necessarily a domain of holomorphy but $`D/\mathrm{\Gamma }`$ can be compactified by adding a finite number of points we can show that the answer to (ii) is affirmative. We do not assume $`D/\mathrm{\Gamma }`$ quasi–projective.
###### Proposition 5.3
Let $`D^n`$ be an open set having a properly discontinous group $`\mathrm{\Gamma }\mathrm{Aut}D`$ such that there exists a compact complex space $`Y`$ with $`D/\mathrm{\Gamma }\mathrm{Reg}Y`$ and $`D/\mathrm{\Gamma }=YS`$, where $`S`$ is a finite set. Then the volume of $`D/\mathrm{\Gamma }`$ in the induced Bergman metric is finite.
###### Demonstration Proof
Since $`M=D/\mathrm{\Gamma }`$ is hyper 1–concave and possesses a positive canonical bundle, we may apply Theorem 3.1 for $`\mathrm{\Gamma }`$ trivial and $`E=K_M`$. As Remark 3.1 (c) shows this gives an upper bound for $`\mathrm{vol}(M)=_M\omega _{}^n/n!`$ . ∎
###### Remark Remark 5.1
We can prove a complete generalization of the asymptotic Morse inequalities of Demailly \[De1\] for the $`L^2`$ cohomology of the covering of a compact manifold $`X`$. For this purpose we elabotate the proof of Theorem 2.1. As there we exploit the idea of Witten–Demailly of constructing a family of subcomplexes of the $`L^2`$–Dolbeault complex having the same cohomology. First let us introduce cohomology. Let us denote by $`N^q(\overline{})`$ the kernel and by $`R^q(\overline{})`$ the range of $`\overline{}`$ , by $`N^q(\overline{}^{})`$ the kernel of $`\overline{}^{}`$ and by $`N^q(k^1\stackrel{~}{\mathrm{\Delta }}_k^{\prime \prime })`$ the kernel of $`k^1\stackrel{~}{\mathrm{\Delta }}_k^{\prime \prime }`$, all acting on $`L_{0,q}^2(\stackrel{~}{X},\stackrel{~}{E}^k\stackrel{~}{F})`$ where $`\stackrel{~}{F}`$ is a $`\mathrm{\Gamma }`$–invariant holomorphic vector bundle of rank $`r`$. We have $`_{(2)}^{0,q}(\stackrel{~}{X},\stackrel{~}{E}^k\stackrel{~}{F}):=N^q(k^1\stackrel{~}{\mathrm{\Delta }}_k^{\prime \prime })=N^q(\overline{})N^q(\overline{}^{}),`$ where the first equality is the definition of the space of harmonic forms and the second is a consequence of the completeness of the metric. If $`q=0`$ then $`_{(2)}^{0,0}(\stackrel{~}{X},\stackrel{~}{E}^k\stackrel{~}{F})`$ coincides to the space of holomorphic $`L^2`$ sections of $`\stackrel{~}{E}^k\stackrel{~}{F}`$. We note also the orthogonal decomposition $`N^q(\overline{})=N^q(k^1\stackrel{~}{\mathrm{\Delta }}_k^{\prime \prime })\overline{R^{q1}(\overline{})}`$ so that
$$_{(2)}^{0,q}(\stackrel{~}{X},\stackrel{~}{E}^k\stackrel{~}{F})=N^q(\overline{})/\overline{R^{q1}(\overline{})}=:H_{(2)}^{0,q}(\stackrel{~}{X},\stackrel{~}{E}^k\stackrel{~}{F})$$
the last group being the (reduced) $`L^2`$ cohomology.
We apply the results of §1 in the following form. Since $`X`$ is compact we can take the set $`\mathrm{\Omega }X`$ to be $`X`$ so that $`\stackrel{~}{\mathrm{\Omega }}=\stackrel{~}{X}`$. We do not use any special metric but take an arbitrary metric on $`X`$ and its pull–back on $`\stackrel{~}{X}`$. Moreover we have $`\stackrel{~}{\mathrm{\Delta }}_k^{\prime \prime }=\stackrel{~}{\mathrm{\Delta }}_k^{\prime \prime }_{\stackrel{~}{\mathrm{\Omega }}}`$. Since $`k^1\stackrel{~}{\mathrm{\Delta }}_k^{\prime \prime }`$ commutes with $`\overline{}`$ it follows that the spectral projections of $`k^1\stackrel{~}{\mathrm{\Delta }}_k^{\prime \prime }`$ commute with $`\overline{}`$ too, showing thus $`\overline{}L_k^q(\lambda )L_k^{q+1}(\lambda )`$ and therefore we have a complex of $`\mathrm{\Gamma }`$–modules of finite $`\mathrm{\Gamma }`$–dimension:
$$0\mathrm{@}>>>L_k^0(\lambda )\mathrm{@}>\overline{}_\lambda >>L_k^1(\lambda )\mathrm{@}>\overline{}_\lambda >>\mathrm{}\mathrm{@}>\overline{}_\lambda >>L_k^n(\lambda )\mathrm{@}>>>0.$$
$`(5.1)`$
$`k^1\stackrel{~}{\mathrm{\Delta }}_k^{\prime \prime }`$ commutes also with $`\overline{}^{}`$ and $`\left(\overline{}_\lambda \right)^{}`$ equals the restriction of $`\overline{}^{}`$ to $`L_k^q(\lambda )`$. Keeping this in mind it is easy to see that
$$N^q(\overline{}_\lambda )/\overline{R^{q1}(\overline{}_\lambda )}=\{uL_k^q(\lambda ):\overline{}_\lambda u=0,\left(\overline{}_\lambda \right)^{}u=0\}=_{(2)}^{0,q}(\stackrel{~}{X},\stackrel{~}{E}^k\stackrel{~}{F}).$$
$`(5.2)`$
We can now apply the following lemma (see \[Sh\]).
###### Algebraic Lemma
Let $`0L_0\mathrm{@}>d_0>>L_1\mathrm{@}>d_1>>\mathrm{}\mathrm{@}>d_n>>L_n0`$ be a complex of $`\mathrm{\Gamma }`$–modules ($`d_q`$ commutes with the action of $`\mathrm{\Gamma }`$ and $`d_{q+1}d_q=0`$). If $`l_q=dim_\mathrm{\Gamma }L_q`$ is finite and $`h_q=dim_\mathrm{\Gamma }H^q(L)`$ where $`H_q(L)=N(d_q)/\overline{R(d_{q1})}`$,
$$\underset{j=0}{\overset{q}{}}(1)^{qj}h_j\underset{j=0}{\overset{q}{}}(1)^{qj}l_j$$
for every $`q=0,1,\mathrm{},n`$ and for $`q=n`$ the inequality becomes equality.
The Algebraic Lemma for the complex (5.1) and relation (5.2) yield
$$\underset{j=0}{\overset{q}{}}(1)^{qj}dim_\mathrm{\Gamma }H_{(2)}^{0,j}(\stackrel{~}{X},\stackrel{~}{E}^k\stackrel{~}{F})\underset{j=0}{\overset{q}{}}(1)^{qj}N_\mathrm{\Gamma }^j(\lambda ,\frac{1}{k}\stackrel{~}{\mathrm{\Delta }}_k^{\prime \prime })$$
for $`q=0,1,\mathrm{},n`$ and for $`q=n`$ the inequality becomes equality. We apply now (1.11):
$$\begin{array}{cc}\hfill \underset{j=0}{\overset{q}{}}(1)^{qj}dim_\mathrm{\Gamma }H_{(2)}^{0,j}(\stackrel{~}{X},\stackrel{~}{E}^k\stackrel{~}{F})k^n(I^q(U,\lambda )& I^{q1}(U,\lambda )+\mathrm{}\hfill \\ & +(1)^qI^0(U,\lambda ))+o(k^n),\hfill \end{array}$$
for $`k\mathrm{}`$. We can now let $`\lambda `$ go to zero through values $`\lambda 𝒩`$. We have thus proved the following.
###### Theorem 5.4
Let $`\stackrel{~}{X}`$ be a Galois covering of group $`\mathrm{\Gamma }`$ of a compact manifold $`X`$. As $`k\mathrm{}`$, the following strong Morse inequalities hold for every $`q=0,1,\mathrm{},n`$ rm :
$$\underset{j=0}{\overset{q}{}}(1)^{qj}dim_\mathrm{\Gamma }H_{(2)}^{0,j}(\stackrel{~}{X},\stackrel{~}{E}^k\stackrel{~}{F})r\frac{k^n}{n!}_{X(q)}(1)^q\left(\frac{ı}{2\pi }𝚌(E)\right)^n+o(k^n).$$
with equality for $`q=n`$ (asymptotic $`L^2`$ Riemann-Roch formula).
In particular $`dim_\mathrm{\Gamma }H_{(2)}^0(\stackrel{~}{X},\stackrel{~}{E}^k\stackrel{~}{F})\frac{k^n}{n!}_{X(1)}\left(\frac{ı}{2\pi }𝚌(E)\right)^n+o(k^n)`$. It follows that if $`E`$ satisfies (D) then for $`k\mathrm{}`$
$$\begin{array}{cc}\hfill dim_\mathrm{\Gamma }H_{(2)}^0(\stackrel{~}{X},\stackrel{~}{E}^k\stackrel{~}{F})& k^n,\hfill \\ \hfill dim_\mathrm{\Gamma }H_{(2)}^q(\stackrel{~}{X},\stackrel{~}{E}^kF)& =o(k^n),q1.\hfill \end{array}$$
Hence the usual dimension of the space of holomorphic $`L^2`$ sections has the same cardinal as $`|\mathrm{\Gamma }|`$ for large $`k`$. This is a generalization of the result of Napier \[Nap\] that $`\stackrel{~}{X}`$ is holomorphically convex with respect to $`\stackrel{~}{E}^k`$ for large $`k`$ if $`X`$ is projective and $`E`$ is positive. If the canonical bundle $`K_X`$ satisfies condition (D), i.e. if there exists a metric $`\omega `$ on $`M`$ such that $`_{X(1)}(\mathrm{Ric}\omega )^n>0`$ where $`X(1)`$ is the set of points where $`\mathrm{Ric}\omega `$ is nondegenerate and has at most one negative eigenvalue, then $`dim_\mathrm{\Gamma }H_{(2)}^0(\stackrel{~}{X},K_{\stackrel{~}{X}}^k)k^n`$ .
###### Remark Remark 5.2
In Proposition 4.1 we have treated the case of a singular hermitian line bundle $`(E,h)`$ over a compact manifold $`X`$. The condition on the singularities were that they are concentrated on an analytic set and moreover the curvature is positive near this analytic set. Then we can work on the complement of the analytic set and by means of the basic estimate study its coverings. If we are interested only in the coverings of $`X`$ then we can rule out the condition of positivity near the singularities. Namely, when the singularities of the metric are algebraic (cf. \[De2\]), Bonavero \[Bon\] shows that the Morse inequlities are true for the cohomology of $`E^k`$ twisted with the corresponding sequence of Nadel’s multiplier ideal sheaves. Given a Galois covering as above we can adapt his proof to estimate the von Neuman dimension of the space $`H_{(2)}^0(\stackrel{~}{X},\stackrel{~}{E}^k_k(\stackrel{~}{h}))`$ of $`L^2`$ holomorphic sections in $`\stackrel{~}{E}^k`$ twisted with the Nadel’s multiplier ideal sheaf coming from the singularities of the $`\mathrm{\Gamma }`$–invariant metric $`\stackrel{~}{h}`$ on $`\stackrel{~}{E}^k`$ (which is the pull–back of a Nadel multiplier ideal sheaf on $`X`$). The conclusion is that when (D) is true, the integral being taken over the regular set of the curvature current, then the von Neuman dimension of $`H_{(2)}^0(\stackrel{~}{X},\stackrel{~}{E}^k_k(\stackrel{~}{h}))`$ grows as $`k^n`$ for large $`k`$.
###### Remark Remark 5.3
Using the approach of this section we can study the growth of the cohomology groups of coverings of $`q`$–convex and $`q`$–concave manifolds. We can either use complete metrics or follow \[GHS\] and use the $`\overline{}`$–Neumann problem setting. Let us give the statements in the latter set-up. Consider a $`q`$–convex manifold $`X`$ in the sense of \[AG\], i.e. there exists a smooth exhausting function $`\phi :X`$ such that $`ı\overline{}\phi `$ has at least $`nq+1`$ positive eigenvalues outside a compact set $`K`$ ($`n=dimX`$ , $`1qn1`$). Consider $`X_c=\{\phi <c\}K`$ with smooth boundary. Then the Levi form of $`bX_c`$ has at least $`nq`$ positive eigenvalues. Let us consider a Galois covering $`\stackrel{~}{X}_d`$ of a bigger sublevel set $`X_dX_c`$ and denote by $`\stackrel{~}{X}_c`$ the induced covering of $`X_c`$. As usual we denote by $`\mathrm{\Gamma }`$ the group of deck transformations. Let us consider also a line bundle $`E`$ over $`X`$ and denote by $`\stackrel{~}{E}`$ its lifting to $`\stackrel{~}{X}_d`$. Both $`\stackrel{~}{X}_c`$ and $`\stackrel{~}{E}`$ come with the liftings of metrics defined on $`X_d`$. We define the (reduced) $`L^2`$ cohomology groups $`H_{(2)}^j(\stackrel{~}{X}_c,\stackrel{~}{E}^k)`$ with respect to these metrics. By \[GHS\] we know that $`dim_\mathrm{\Gamma }H_{(2)}^j(\stackrel{~}{X}_c,\stackrel{~}{E}^k)<\mathrm{}`$ for $`jq`$. With the method used in this paper we can prove that for $`jq`$ and $`k\mathrm{}`$ :
The proof consists of showing that the basic estimate holds in bidegree $`(0,j)`$ on $`\stackrel{~}{X}_c\stackrel{~}{L}`$, where $`L`$ is a compact set of $`X_c`$ , for forms satisfying the $`\overline{}`$–Neumann conditions on $`b\stackrel{~}{X}_c`$ . This is achieved using the liftings of the metrics constructed in \[AV\] where the case $`\mathrm{\Gamma }`$ trivial is treated. Then we can apply again the analysis from §1. If $`E`$ is $`q`$–positive outside $`K`$ then the leading term in (1) simplifies as shown in (2). These estimates were obtained in the case $`\mathrm{\Gamma }=\{\mathrm{Id}\}`$ in \[Bou\] for certain complete metrics on $`X_c`$ which permit to prove the same inequalities for the full cohomology group $`H^j(X_c,E^k)`$. For the case of coverings we have to restrict ourselves to $`L^2`$ cohomology groups. As for coverings of $`q`$–concave manifolds we get the same conclusion as in (1) for $`jnq1`$. The nice simplification of the leading term holds if we impose a negativity condition outside a compact set. However there are cases of concave manifolds and positive bundles for which we have an effective estimate of $`dim_\mathrm{\Gamma }H_{(2)}^0(\stackrel{~}{X}_c,\stackrel{~}{E}^k)`$, see §3.
## §6 Weak Lefschetz theorems
Nori \[No\] generalized the Lefschetz hypersurface theorem. Assume $`X`$ and $`Y`$ are smooth connected projective manifolds and $`Y`$ is a hypersurface in $`X`$ with positive normal bundle and $`dimY1`$. Then the image of $`\pi _1(Y)`$ in $`\pi _1(X)`$ is of finite index. Recently, Napier and Ramachandran \[NR\] proposed an analytic approach and generalized Nori’s theorem showing that $`Y`$ may have arbitrary codimension (but $`dimY1`$). They use the $`\overline{}`$–method on complete Kähler manifolds to separate the sheets of appropriate coverings. In the sequel we use the Riemann–Roch inequalities to study non–necessarily Kähler manifolds. However our method requires that the image group is normal since we can deal only with Galois coverings. First we introduce the notion of formal completion. Let $`Y`$ be a complex analytic subspace of the manifold $`U`$ and denote by $`_Y`$ the ideal sheaf of $`Y`$. The formal completion $`\widehat{U}`$ of $`U`$ with respect to $`Y`$ is the ringed space $`(\widehat{U},𝒪_{\widehat{U}})=(Y,proj\; lim𝒪_U/_Y^\nu )`$. If $``$ is an analytic sheaf on $`U`$ we denote by $`\widehat{}`$ the sheaf $`\widehat{}=proj\; lim(𝒪/_Y^\nu )`$. If $``$ is coherent then $`\widehat{}`$ is too. Moreover by Proposition VI.2.7 of \[BS\] the kernel of the mapping $`H^0(U,)H^0(\widehat{U},\widehat{})`$ consists of the sections of $``$ which vanish on a neighbourhood of $`Y`$. Hence for locally free $``$ the map is injective.
###### Theorem 6.1
Let $`M`$ be a hyper $`1`$–concave manifold carrying a line bundle $`E`$ which satisfies (D) and is semi-positive outside a compact set. Let $`Y`$ be a connected compact complex subspace of $`M`$ satisfying: (i) for any $`k`$, $`dimH^0(\widehat{M},\widehat{}_k)<\mathrm{}`$, where $`_k=𝒪(E^kK_M)`$, (ii) the image $`G`$ of $`\pi _1(Y)`$ in $`\pi _1(X)`$ is normal in $`\pi _1(X)`$. Then $`G`$ is of finite index in $`\pi _1(X)`$.
###### Demonstration Proof
We follow the proof given in \[NR\]. Since $`G`$ is normal there exists a connected Galois covering $`\pi :\stackrel{~}{M}M`$ such that the group of deck transformations is $`\mathrm{\Gamma }=\pi _1(M)/G`$. The cardinal $`|\mathrm{\Gamma }|`$ equals the index of $`G`$ in $`\pi _1(M)`$. Let $`\stackrel{~}{E}=\pi ^1E`$. By Theorem 3.1, there exists $`C>0`$ such that for large $`k`$, $`dim_\mathrm{\Gamma }H_{(2)}^{n,0}(\stackrel{~}{M},\stackrel{~}{E}^k)Ck^n`$. Let us choose a small open neighbourhood $`V`$ of $`Y`$ such that $`\pi _1(Y)\pi _1(V)`$ is an isomorphism; so the image of $`\pi _1(V)`$ in $`\pi _1(M)`$ is $`G`$. Hence, if we denote by $`ȷ`$ the inclusion of $`V`$ in $`M`$, there exists a holomorphic lifting $`\stackrel{~}{ȷ}:V\stackrel{~}{M}`$, $`\pi \stackrel{~}{ȷ}=ȷ`$. Since $`\stackrel{~}{ȷ}`$ is locally biholomorphic the pull–back map $`\stackrel{~}{ȷ}^{}:H_{(2)}^{n,0}(\stackrel{~}{M},\stackrel{~}{E}^k)H^{n,0}(V,E^k)`$ is injective. On the other hand $`H^0(V,_k)H^0(\widehat{V},\widehat{}_k)=H^0(\widehat{M},\widehat{}_k)`$. By (i) the latter space is finite dimensional so $`dimH_{(2)}^{n,0}(\stackrel{~}{M},\stackrel{~}{E}^k)<\mathrm{}`$. We know that $`dim_\mathrm{\Gamma }H_{(2)}^0(\stackrel{~}{M},\stackrel{~}{E}^kK_{\stackrel{~}{M}})>0`$ for $`k>C^{1/n}`$. If $`\mathrm{\Gamma }`$ were infinite this would yield $`dimH_{(2)}^{n,0}(\stackrel{~}{M},\stackrel{~}{E}^k)=\mathrm{}`$ which is a contradiction. Therefore $`|\mathrm{\Gamma }|<\mathrm{}`$ and $`dimH_{(2)}^{n,0}(\stackrel{~}{M},\stackrel{~}{E}^k)C|\mathrm{\Gamma }|k^n|\mathrm{\Gamma }|`$ for $`k>C^{1/n}`$. Thus $`|\mathrm{\Gamma }|dimH^0(\widehat{M},\widehat{}_k)`$ for large $`k`$. ∎
###### Remark Remark 6.2
(a) By a theorem of Grothendieck \[Gro\], condition (i) is fulfilled if $`Y`$ is locally a complete intersection with ample normal bundle $`N_Y`$ (or $`k`$–ample in the sense of Sommese, $`k=dimY1`$).
(b) We can replace condition (i) with the requirement that $`Y`$ has a fundamental system of pseudoconcave neighbourhoods $`\{V\}`$. Then $`dimH^0(V,_k)`$ is finite by \[An\]. This happens for example if $`Y`$ is a smooth hypersurface and $`N_Y`$ has at least one positive eigenvalue or, if $`Y`$ has arbitrary codimension, if $`N_Y`$ is sufficiently positive in the sense of Griffiths \[Gri1\].
(c) Condition (ii) is trivially satisfied if $`\pi _1(Y)=0`$. Thus, if $`M`$ contains a simply connected subvariety satisfying either (a) or (b), $`\pi _1(M)`$ is finite.
(d) By Corollary 3.6, Theorem 6.1 can also be applied to the pertubed structures considered there.
Using Proposition 4.3 we can can show that Nori’s theorem holds for all Moishezon spaces $`X`$.
###### Theorem 6.2
Let $`X`$ be an irreducible reduced normal Moishezon compact complex space and let $`E`$ be the (positive in the sense of currents) line bundle given by Theorem 4.2. Suppose that $`M`$ is a Zariski open set of $`X`$ and $`Y\mathrm{Reg}(M)`$ be a connected compact complex subspace such that for any $`k`$, $`dimH^0(\widehat{M},\widehat{}_k)<\mathrm{}`$, where $`_k=𝒪(E^k)`$. Then the image $`G`$ of $`\pi _1(Y)`$ in $`\pi _1(M)`$ is of finite index in $`\pi _1(M)`$.
###### Demonstration Proof
Since $`X`$ is normal we have an isomorphism $`\pi _1(\mathrm{Reg}M)\pi _1(M)`$, so that we may assume $`M\mathrm{Reg}(X)`$. We find a connected unramified covering $`p:\stackrel{~}{M}M`$ such that $`p_{}\pi _1(\stackrel{~}{M})=G`$. If $`d`$ is the number of sheets, $`d=|\pi _1(M)/G|`$, the index of $`G`$ in $`\pi _1(M)`$. The preceding proof applies by using Proposition 4.3 instead of Theorem 3.1. and the usual dimension instead of the $`\mathrm{\Gamma }`$–dimension. ∎
Note that Napier and Ramachandran also considered cases when $`X`$ is not necessarily projective, but their result does not imply diectly Theorem 6.2. |
warning/0002/astro-ph0002255.html | ar5iv | text | # Multiwavelength observations of Mkn 501 during the 1997 high state
## 1 Introduction
The BL Lac object Mkn 501 is very close ($`z=0.0337`$, Ulrich et al. ulrich (1975)) and has been studied extensively at all wavelengths. Together with its sister object Mkn 421, it was among the two first BL Lac objects with known radio (Colla et al. colla (1972)), X-ray (Schwartz et al. schwartz (1978)) and TeV gamma-ray (Quinn et al. quinn (1996), Bradbury et al. bradbury (1997)) counterparts. Recently, it was also marginally detected at photon energies $`>100`$ MeV by the EGRET instrument on board CGRO (Kataoka et al. kataoka (1999)).
During 1997, the object was found to be in an extreme high state with a TeV flux on average 20 times higher than in 1996 (Breslin et al. breslin (1997)). The source exhibited strong variability on timescales of days with a possible quasi-periodically varying component with a timescale of about 25 days (Kranich et al. kranich (1999)). To complete the list of reasons for excitement for the observers, Mkn 501 reached in some of its flares fluxes of more than 10<sup>-10</sup> cm<sup>-2</sup>s<sup>-1</sup> (above 1.5 TeV) - the most intense TeV emission ever measured so far from any astronomical object. However, the shortest observed variability timescale (5 hours, Aharonian et al. 1999a ) was significantly longer than that observed for Mkn 421 (Gaidos et al. gaidos (1996)).
All available TeV observatories monitored the event for several months (Samuelson et al. whipple (1998), Aharonian et al. 1999a ,1999b , Hayashida et al. telarray (1998), Djannati-Ataï et al. cat (1999)). The most complete dataset was produced by the HEGRA Cherenkov Telescope “1” (CT1) (Aharonian et al., 1999b ) which was even able to observe Mkn 501 under the presence of moonlight, however with reduced sensitivity, thereby filling many gaps in the lightcurve. This telescope also obtained the confirmatory observations in 1996.
Also from HEGRA come probably the most accurate spectral measurements in the TeV regime. They were carried out by the HEGRA system of (at the time 4) Cherenkov telescopes (CTS) and are largely concurrent with the CT1 measurements, however with less time coverage (Aharonian et al. 1999a ).
The origin of the TeV $`\gamma `$-ray emission and the reasons for its variability are still essentially unknown. The most popular models explain the TeV emission as near-infrared to UV photons which have been upscattered via the inverse Compton effect by very high energy electrons which are known from radio, optical and X-ray observations to be present in the jets of BL Lac objects. Possible sources of the seed photons for Compton scattering are the synchrotron radiation produced within the jet by the same population of relativistic electrons (synchrotron self-Compton model, SSC; Marscher & Gear marscher (1985); Maraschi, Celotti & Ghisellini maraschi (1992); Bloom & Marscher bloom (1996)), or radiation from outside the jet (external inverse Compton model, EIC). This external radiation could be the quasi-thermal radiation field of an accretion disk surrounding a supermassive black hole which is generally believed to power the relativistic jets. The accretion disk radiation can enter the jet either directly (Dermer, Schlickeiser & Mastichiadis dermer92 (1992), Dermer & Schlickeiser dermer93 (1993)) or after being rescattered by circumnuclear material (Sikora, Begelman & Rees sikora (1994), Blandford & Levinson blandford (1995), Dermer, Sturner & Schlickeiser dermer97 (1997)). It is also possible that synchrotron radiation produced within the jet and reflected by circumnuclear debris is the dominant source of soft photons during flares (Ghisellini & Madau ghisellini96 (1996), Bednarek bednarek (1998)), although it has been shown that this process is unlikely to be efficient in the case of BL Lac objects (Böttcher & Dermer boettcher98 (1998)).
As an alternative, Mannheim (mannheim93 (1993)) has suggested a hadronic model in which protons are the primarily accelerated particles in the jet and the $`\gamma `$-ray emission is produced by secondary pions and electron-positron pairs produced in photopion and photopair production interactions of the ultrarelativistic protons in the jet with external radiation. The time-averaged broadband spectrum of Mkn 501 has been fitted using this model by Mannheim et al. (mannheim96 (1996), mannheim98 (1998)). However, the attempts to explain the short and intermediate-term variability of blazars with hadronic jet models have only just started (Rachen & Mannheim rachen99 (1999)). For this reason, we concentrate on leptonic jet models in this paper since the variability time scales predicted by these models are in good agreement with the observed intraday variability of blazars.
Recently, Fossati et al. (fossati (1997)) have compared the broadband spectra of different types of $`\gamma `$-ray emitting AGN and suggested a continuous sequence FSRQs (flat spectrum radio quasars) $``$ LBLs (low-frequency peaked BL Lac objects) $``$ HBLs (high-frequency peaked BL Lac objects), characterized by decreasing bolometric luminosity, decreasing dominance of the total energy output in $`\gamma `$-rays compared to the emission at lower frequencies and a shift of the peak frequencies of the synchrotron and the $`\gamma `$-ray component towards higher frequencies. Recent modelling efforts of various blazar-type AGN have revealed that this sequence is consistent with a decreasing importance of external radiation as a source of soft photons for Compton scattering in the jet (Ghisellini et al. ghisellini98 (1998)). This suggests that the extreme HBLs like Mkn 501 or Mkn 421 can be well fitted with strongly SSC-dominated jet models, as was shown, e. g., by Mastichiadis & Kirk (mastichiadis97 (1997)) and Pian et al. (pian (1998)).
From the SSC model one expects a strong correlation between the X-ray and the TeV emission. And indeed, by comparing the daily TeV measurements with daily averages from the All Sky Monitor (ASM) on board the Rossi X-Ray Timing Explorer (RXTE), a significant correlation with a most probable time-lag of 0 days (no time-lag) was found (Aharonian et al. 1999a ; 1999b ). Furthermore, a clear overall X-ray high state was visible in the 1997 ASM measurements which coincided with the TeV high state (see also figure 1). And, as discovered by BeppoSAX, the X-ray peak of the spectral energy distibution (SED) had shifted from 10-20 keV in 1996 to 100-200 keV in 1997 (Pian et al. pian (1998)), consistent with the assumption that the X-ray and TeV-$`\gamma `$-ray flares are produced by a more powerful electron acceleration, shifting the high-energy cutoff of the electron distribution to higher energies and hardening the X-ray spectrum. This observation was, however, only made during one flare in 1997. In this paper we give evidence that the hardening took place during all flares in 1997. An alternative explanation for the synchrotron spectral changes, in terms of a steadily-emitting helical-jet model, has been presented by Villata & Raiteri (villata99 (1998)).
Apart from the short-term variability on timescales of hours to days, there is also a longer-term variability in Blazars which has so far been investigated mainly in the optical regime (see e.g. Katajainen et al. (katajainen (1999)) and references therein) and recently also in X-rays (e.g. Mc Hardy mchardy (1999)). This variability shows remarkable amplitudes (e.g. 4.7 mag for Mkn 421 in the optical) and in the case of OJ 287 there is even evidence for periodicity.
In this respect Mkn 501 is not yet very well explored. Mkn 501 has been observed in the TeV regime since its discovery in 1995. The source showed low emission close to the sensitivity limits until the onset of the 1997 high state. Since the duty cycle of the TeV observatories is only about 10 %, this does not prove the absence of strong short flares prior to this high state, but the increased average intensity of the source can be described as an increase in flaring probability from 1996 to 1997 by at least an order of magnitude. This description is especially appropriate since even during the high state, the source returns to quiescent (comparable with 1996) levels of emission for periods of up to a few days. These transitions are seen irrespective of the presence of short, strong flares on time-scales of a day.
The last observations in 1996 were made in August and found Mkn 501 still quiescent while the first observations in 1997 were made in February and found the source already flaring. The transition from low to high flaring probability obviously took place on timescales shorter than half a year. The correlated X-ray data from the RXTE ASM confirm this and show in addition that the change has been a smooth process over several months (see e.g. figure 8 in Aharonian et al. 1999b ).
In this article we explore the multi-wavelength variability of Mkn 501 over medium timescales, and try to relate its behaviour to the physical parameters of a leptonic jet model. For this purpose we construct weekly SEDs using the HEGRA CT1 flux data and HEGRA CT System spectral data together with data from longer wavelengths, namely radio, optical, soft X-ray and hard X-rays. See table 1 for an enumeration of the instruments and energy ranges.
This paper may represent an important step toward understanding the origin of strong changes in the flaring probability of Mkn 501 and Blazars in general, but our data is clearly not sufficient to give a final answer to this question.
All data except that from HEGRA and the RXTE ASM are published here for the first time. The BATSE data are especially valuable since they confine the intensity at the X-ray peak of the SED.
## 2 Observations and Data Analysis
The observations used to fit the multi-wavelength spectrum cover mostly the synchrotron part of the SED. Only the TeV data explore what is believed to be the Inverse Compton emission, although OSSE and EGRET observations cover a few days in 1996 and 1997. As a guide line for our fitting procedure we take into account the highest ever observed EGRET flux from Mkn 501, $`F(>100\mathrm{MeV})=(32\pm 13)10^8`$ photons cm<sup>-2</sup> s<sup>-1</sup> with a photon spectral index $`\alpha =1.6\pm 0.5`$, measured during 1996 Mar 25 - 28 (Kataoka et al. kataoka (1999)), as an upper limit.
In the definition of time-bins for the multi-wavelength dataset, we start by subdividing the HEGRA CT1 lightcurve and bin all other data accordingly.
The data from HEGRA are binned in time from March through October 1997 (MJD 50514 - MJD 50708) resulting in 28 equidistant time bins. This weekly temporal resolution is an order of magnitude larger that the longest observed TeV intraday variability timescale of 15 h (Aharonian et al. 1999a ) and is believed appropriate, given the nature of the available data and the medium scale variability timescale which we have chosen to explore.
However, the HEGRA points are not spread as uniformly over time as are the BATSE and RXTE data. The “center of gravity in observation time” (defined as the weighted mean of the observation time of the daily points each weighted by the duration of the individual observation) from HEGRA would therefore in general not coincide with that from BATSE and RXTE. In order to compensate for this, we calculate the center of gravity in observation-time for each weekly HEGRA timebin and use these as time-bin-centers for the other data. The edges of these bins are then defined by the average of two adjacent bin-centers.
Figure 1 shows the lightcurves from each instrument which went into this analysis. The weekly points after the binning are given in table 2. The first column in the table gives the bin-centers as described above. The following subsections describe the data used to compile the table.
### 2.1 Optical and Radio Observations
The optical observations were performed using the telescopes and filters listed in table 1. All the observations were made with CCD-cameras. All CCD-images were treated the normal way with flat field and bias corrections. The magnitudes were measured using either DAOPHOT in IRAF (NOT and Tuorla data; for more details see Katajainen et al. katajainen (1999)), or with the ROBIN-procedure developed in Torino (e.g. Villata et al. villata97 (1997)), or using the automated reduction routine developed in Perugia (e.g. Tosti et al. tosti (1996)). All magnitudes were measured using 10 arcsecond aperture. The use of the same aperture is important because Mkn 501 has a large host galaxy (e.g. Nilsson et al. nilsson (1999)). For the calibration we used the calibration sequence given by Fiorucci & Tosti (fiorucci (1996)) and Villata et al. (villata98 (1998)).
For some of the weekly observing periods, no optical data were available. However, the optical flux is extremely important for constraining the spectral index of the synchrotron spectra in the optical – X-ray range, which, in turn, is essential for constraining the spectral index of the electron spectrum. As can be seen in table 2, the optical flux exhibits only moderate variability on the time scales considered here. Thus, for fitting purposes we assume that the optical flux during those viewing periods for which no optical observations were available, was within the range of optical fluxes observed during the whole campaign which yields $`\nu F_\nu (opt.)=(7.485\pm 0.774)10^{11}`$ erg cm<sup>-2</sup> s<sup>-1</sup>.
The radio observations were made with the Metsähovi Radio Telescope at 22 GHz as part of the ongoing quasar monitoring program, which started in 1980. Currently about 85 sources, mainly Northern flat spectrum quasars, should be observed monthly at 22 and 37 GHz. The total number of observations is now over 40000. For details of the observing strategy and reductions, as well as the data until 1995.5 see Teräsranta et al. (terasranta (1998)).
The integration time of each point was essentially the same, so single points were formed from the data points in the individual timebins by simple averaging.
### 2.2 Soft X-Ray Data
Since the beginning of 1996, the all-sky monitor (ASM) of the Rossi X-ray Timing Explorer (RXTE) satellite has been observing Mkn 501 in the 2-12 keV energy band. The data used in this analysis are taken from the publicly available ASM data products provided by the ASM/RXTE teams at MIT and at the RXTE SOF and GOF at NASA’s GSFC. These measurements are given by the authors (see e.g. Levine et al. levine (1996)) as rates $`R_1,R_2,R_3`$ in units of counts per second in three energy bins: 1.3-3.0 keV, 3.0-5.0 keV, and 5.0-12.1 keV. From these we calculate the total rate $`R`$ by summing the three bins.
For Mkn 501 we obtain by averaging over the period MJD 50510 - 50710
$$R_{mkn501}=1.3\pm 0.4$$
(1)
In order to assess the X-ray spectrum and flux of Mkn 501, we use the ASM data for the Crab Nebula which is publicly available from the same source. The average rate of the Crab Nebula during the period MJD 50510 – 50710 is
$$R_{crab}=75.7\pm 0.5$$
(2)
Hence, in soft X-rays, Mkn 501 is even during this high state a significantly weaker source than the Crab Nebula.
In the energy range 2 keV to 60 keV, the spectrum of the Crab Nebula is a stable power law. The spectral index of the differential photon spectrum is known with an accuracy of 1% to be $`\alpha _{crab}=2.1`$ (e.g. Toor & Seward toor (1974), Pravdo & Serlemitsos pravdo (1981) or Pravdo, Angelini & Harding pravdo-b (1997)). The emission measured by the RXTE ASM is the sum of the steady Nebula and the pulsed Crab Pulsar emission. The latter has on average a harder spectrum than the nebula. Below 12.1 keV, however, the pulsed fraction of this emission is $`<8`$ %. We can therefore use the knowledge of the Crab Nebula spectrum to derive an approximate relative calibration factor $`k`$ for the rates $`R_2`$ and $`R_3`$ using
$$\frac{kR_{2,crab}}{R_{3,crab}}=\frac{E_3^{1\alpha _{crab}}E_2^{1\alpha _{crab}}}{E_4^{1\alpha _{crab}}E_3^{1\alpha _{crab}}}=1.213$$
(3)
where $`E_2=3.0`$ keV, $`E_3=5.0`$ keV, and $`E_4=12.1`$ keV are the energy bin edges of the second and third energy bin. We ignore the first energy bin since it is strongly influenced by interstellar X-ray absorption which is dependent on the column density and hence varies between sources (Remillard remillard (1999)).
Averaging over the available Crab data from the observation period under discussion (MJD 50510 - 50710) we obtain $`k=1.318\pm 0.0024`$. There is no indication of a variability of the value of $`k`$ (see figure 2). Hence we can assume that it is also valid for the Mkn 501 observations of the same detector. In order take into account that there is a pulsed component with a harder spectrum, we add an additional error of 8 % to the error of $`R_3`$ and obtain thus
$$k=1.32\pm 0.023$$
(4)
For Mkn 501 we calculate the spectral index $`\alpha `$ of the differential photon spectrum using the equation
$$\frac{kR_{2,mkn501}}{R_{3,mkn501}}=\frac{E_3^{1\alpha }E_2^{1\alpha }}{E_4^{1\alpha }E_3^{1\alpha }}$$
(5)
and varying $`\alpha `$ until the two sides of the equation are equal to an accuracy better than 0.1 %. The error $`\delta \alpha `$ of this index, we estimate from the approximate formula
$$\alpha =\frac{\mathrm{log}(kR_2/(5.03.0))\mathrm{log}(R_3/(12.15.0))}{\mathrm{log}((3.0+5.0)/(5.0+12.1))}$$
(6)
which leads to
$$\delta \alpha =\sqrt{1.733((\frac{\delta k}{k})^2+(\frac{\delta R_2}{R_2})^2+(\frac{\delta R_3}{R_3})^2)}$$
(7)
where $`\delta k`$, $`\delta R_2`$, and $`\delta R_3`$ are the errors of the corresponding quantities. The time resolved values for $`\alpha _{mkn501}`$ are shown in figure 3. The same method applied to the Crab Nebula data (however with $`\delta k=0.0024`$) yields, as expected, on average the spectral index we have put in. Also this is shown in figure 3. From a constant fit to the values in this figure, we obtain
$$\alpha _{mkn501}=1.76\pm 0.04$$
(8)
and
$$\alpha _{crab}=2.10\pm 0.02$$
(9)
Thus we find that during the 1997 high state, Mkn 501 had on average a significantly harder soft X-ray spectrum than the Crab Nebula. The spectral variability of Mkn 501 is below the ASM’s sensitivity. The distribution of the weekly values is consistent with a constant value (reduced $`\chi ^2`$ = 0.93) and so is that for the Crab Nebula (reduced $`\chi ^2`$ = 0.95).
In order to calculate the $`\nu F_\nu `$ values in erg cm<sup>-2</sup> s<sup>-1</sup>, we use the knowledge of the flux of the Crab Nebula. Here we face the problem that there is still a disagreement of up to 25 % between the measured normalization constants of the Crab Spectrum from different experiments although there is perfect agreement in the spectral index. The discrepancy seems to stem from not very well understood systematic differences between the detectors (see the discussion in Pravdo, Angelini & Harding pravdo-b (1997)). We use the two most extreme recent measurements of the differential photon flux of the Crab Nebula together with the Crab Pulsar (pulse-averaged) and take their average as our normalization and their difference as the systematic error of this quantity. From Pravdo & Serlemitsos (pravdo (1981)) we get at 5.2 keV a flux of (0.236 $`\pm `$ 0.006) photons cm<sup>-2</sup>s<sup>-1</sup>keV<sup>-1</sup>, while from Pravdo, Angelini & Harding (pravdo-b (1997)) we get (0.302 $`\pm `$ 0.001) photons cm<sup>-2</sup>s<sup>-1</sup>keV<sup>-1</sup>. The average of these values corresponds to a differential energy flux of
$$F_{crab}(5.2\mathrm{keV})=(2.24\pm 0.27_{syst})\times 10^9\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1\mathrm{keV}^1$$
(10)
where the systematic error is the difference between the averaged values divided by 2. The statistical errors of the two measurements are negligible. The energy flux of any other source with ASM rates $`R_2`$ and $`R_3`$ and differential photon spectral index $`\alpha `$ is then calculated by
$$\begin{array}{ccc}\hfill \nu F_\nu [\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1]& =& (R_2+R_3)[\mathrm{s}^1]\frac{F_{crab}\left(5.2\mathrm{keV}\right)\left[\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1\mathrm{keV}^1\right]}{\left(R_{2,crab}+R_{3,crab}\right)\left[\mathrm{s}^1\right]}\hfill \\ & & 5.2^{1+\alpha _{crab}\alpha }\frac{1\alpha }{1\alpha _{crab}}\frac{E_4^{1\alpha _{crab}}E_2^{1\alpha _{crab}}}{E_4^{1\alpha }E_2^{2\alpha }}\hfill \\ & =& R[\mathrm{s}^1]\frac{5.2^{1\alpha }\left(1\alpha \right)}{12.1^{1\alpha }3.0^{1\alpha }}3.13\times 10^{10}\hfill \end{array}$$
(11)
where we correct for the difference in the spectra of Crab Nebula and the source in question by making the Ansatz that the measured differential rates have the same ratio as the differential fluxes. Furthermore we have used the result $`R_{2,crab}+R_{3,crab}=48.5\pm 0.7`$ s<sup>-1</sup> obtained from the dataset under discussion. Inserting the spectral index of Mkn 501 gives:
$$\nu F_\nu (\mathrm{Mkn501},5.2\mathrm{keV})=(R_2+R_3)[\mathrm{s}^1]2.40\times 10^{10}\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1$$
(12)
The error of this flux value is determined by propagating all errors of the quantities involved which gives
$$\begin{array}{cc}\delta (\nu F_\nu )(\mathrm{Mkn501},5.2\mathrm{keV})=\hfill & \\ (\sqrt{5.76\times 10^{20}((\delta R_2)^2+(\delta R_3)^2)+1.21\times 10^{23}(R_2+R_3)^2}+\mathrm{\hspace{0.17em}0.12}\nu F_\nu )\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1\hfill & \end{array}$$
(13)
where the term outside the square-root stems from the systematic error of the Crab flux normalization.
### 2.3 Hard X-Ray Data
The Hard X-ray fluxes were measured using the Burst and Transient Source Experiment (BATSE) onboard the CGRO satellite. Although BATSE is an uncollimated detector, accurate point source fluxes can be measured using the Earth Occultation method described in Harmon et al. (batsemethod (1992)). Daily sensitivities are 100 mCrab and over integrations of years, sources as weak as 3 mCrab can be detected. Several months are usually needed to obtain a statistically significant flux from Mkn 501 with BATSE, but intense daily flares can be seen, and the overall high state of 1997 allowed useful measurements in each weekly interval even outside the flares.
The source is measured only when it sets and rises from behind the Earth. Since two such occultations occur per spacecraft orbit (roughly every 90 minutes), up to 32 independent flux measurements can be made per day. There exists, however, a wide variation in this number because of passage of the spacecraft through the South Atlantic Anomaly, telemetry gaps from loss of TDRSS contact, and other events, which occur randomly relative to the Mkn 501 steps. For the data shown here between 61 and 211 measurements went into an individual weekly point. Each measurement lasts about 8 seconds giving a duty cycle of about 0.2%. In a 7 day time bin during the period discussed in this paper, BATSE observes of the order of $`10^5`$ photons from the source (more during the flares).
The fluxes are integrated between 20 and 200 keV, and are calculated by folding the measured counts through the BATSE detector response assuming a differential source powerlaw spectrum of index -2.0. The -2.0 spectrum is the best fit to the flare measured on MJD 50550-51 (between 20 and 1000 keV). Spectra for other time intervals were also calculated and were consistent with -2.0. Uncertainties of 10 % in the spectral index during flare times, larger outside, make spectral variability difficult to assess for this source, so that the index of -2.0 was used for each weekly interval. Fluxes were calculated for smaller energy bands but the single 20-200 keV (median energy 36.4 keV) point for each interval is the most useful outside intense flares.
Seven of the weekly averages are not statistically significant, and one period shows a $`1.8\sigma `$ deficit. These eight points are inconsistent with a zero-level flux at the 99 % confidence level. The seven low but positive points are inconsistent with a zero-flux level at the 98-99 % confidence level and are well fit by a constant flux of $`1.1\pm 0.3\times 10^{10}`$ erg cm<sup>-2</sup> s<sup>-1</sup>, implying that a flux is present which is below the sensitivity of the BATSE instrument in weekly integrations. An analysis of 30 blank fields on the sky shows that with 100 weekly integrations for each field one can not distinguish (using an F-test) between a zero-flux level and the weighted mean of the weekly averages for any of the 30 fields. This implies that systematic effects are not responsible for the excess seen in the low points and that a flux is indeed present. For this reason these low but positive values have been included as detections in the multiwavelength fits.
### 2.4 TeV data
As described above, we use the data from the CT1 lightcurve of Mkn 501 (integral flux above 1.5 TeV) published by Aharonian et al. (1999b ). These data are available in daily points based on observations of between 0.5 h and 5 h duration each. We group these points according to our weekly time-bins and calculate an average flux from the up to seven values weighting each daily point by its observation time.
Apart from giving daily flux measurements, the HEGRA papers Aharonian et al. (1999a ), (1999b ), and (1999c ) also determine the average spectral shape in the range between 0.5 and $``$ 25 TeV with high accuracy. They find
$$dF/dE=N_0E^\alpha \mathrm{exp}(E/E_0)$$
(14)
where $`N_0=(10.8\pm 0.2\pm 2.1)\times 10^{11}`$ cm<sup>-2</sup> s<sup>-1</sup> TeV<sup>-1</sup>, $`\alpha =1.92\pm 0.03\pm 0.20`$, and $`E_0=(6.2\pm 0.4\pm 2.2)`$ TeV. The first error given is the statistical, the second the systematic error. Furthermore, they find that there is no spectral variability up to their sensitivity of $`\delta \alpha 0.1`$ on all relevant timescales.
In order to include this important spectral information in our model fit, we make the assumption that there is indeed no spectral variability and extrapolate the points measured at 1.5 TeV by CT1 up to 10 TeV and down to 0.8 TeV. From the average integral flux values $`F_{1.5}`$ in units of photons cm<sup>-2</sup>s<sup>-1</sup> we obtain the $`\nu F_\nu `$ values at photon energy $`E`$ in TeV using
$$\nu F_\nu (E)[\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1]=1.6022F_{1.5}\frac{E^{(2\alpha )}\mathrm{exp}(E/E_0)}{\underset{1.5}{\overset{\mathrm{}}{}}𝑑ϵϵ^\alpha \mathrm{exp}(ϵ/E_0)}.$$
(15)
The extrapolation uses the measured spectral shape (equation 14) and fully propagates all statistical errors (error of the CT1 point, error of the spectral index $`\alpha `$ and the error of the cut-off energy $`E_0`$) to form the error of the extrapolated point. The systematic errors of the spectral shape are expected to influence all measured TeV spectra in the same way since they are caused by the 15 % uncertainty in the absolute energy calibration of the Cherenkov telescopes. Thus, for the purpose of a spectral variability study, the statistical errors are those which actually determine the uncertainty in the differences of the spectral shape between different measurements. For this reason we propagate only the statistical errors in the extrapolation. In this way, the columns 6 and 8 of table 2 are obtained. In the model fit (section 2.5) we take into account the uncertainty in the energy scale by introducing errors of $`\pm 15`$ % along the energy axis.
The energies to which we extrapolate are chosen as a compromise of maximum distance to 1.5 TeV and minimum systematic errors. The latter are increasing up to several 10 % towards both ends of the range for which the spectrum has been measured (see Aharonian et al. 1999c , figure 9) but are still small at 0.8 and 10 TeV.
The correlation between the three TeV points which we introduce into the model fit by performing the described extrapolation is not problematic since we do not plan to calculate absolute $`\chi ^2`$ values in the fits for proving that the model describes the data better than another. Instead, the fits serve the aim to study the time-dependent behaviour of the model parameters. The extrapolated points are only a means to take into account the available spectral information.
### 2.5 Fitting the model to the SEDs
To each weekly SED we fit the Blazar jet model described in detail by Böttcher et al. (boettcher97 (1997)). The model assumes that isolated components (blobs) of relativistic pair plasma, which are assumed to be spherical in the co-moving frame, are injected instantaneously into the jet, and follows the self-consistent evolution of the particle and radiation spectra as the blob moves outward along the jet, taking into account all relevant radiation, cooling, and absorption mechanisms: synchrotron radiation, synchrotron self-absorption, synchrotron self-Compton scattering, external Compton scattering of direct accretion disk radiation, $`\gamma `$-$`\gamma `$ absorption and pair production intrinsic to the source. The magnetic field is assumed to decay along the jet as $`Br^1`$, where $`r`$ is the distance from the center of the AGN. The emerging, time-averaged spectra are corrected for $`\gamma \gamma `$ absorption by the intergalactic infrared background radiation using the lower model spectrum given by Malkan & Stecker (malkan98 (1998)).
Since we are interested in weekly averages, the emission from single blobs is time-integrated over the jet evolution and subsequently re-converted into a flux by dividing the fluence by an average repetition time $`\mathrm{\Delta }t_{\mathrm{rep}}`$ of blob ejection events. We assume that a fraction $`f<1`$ of the jet is filled with relativistic pair plasma. The filling factor is given by $`fR_B^{}/(\mathrm{\Gamma }c\mathrm{\Delta }t_{\mathrm{rep}})`$, where $`R_B^{}`$ is the blob radius in the co-moving frame and $`\mathrm{\Gamma }`$ is the bulk Lorentz factor of the blobs. Even if the filling factor is close to unity, it is a reasonable approximation to assume that the blobs do not interact with each other because the synchrotron and SSC radiation produced within the jet are isotropic in the comoving frame so that most of the radiation escapes to the sides without interaction with the rest of the jet.
As mentioned in the introduction, extreme HBLs like Mkn 501 or Mkn 421 are generally well described by a pure SSC model. A simple analytic estimate shows that the observed TeV $`\gamma `$-ray spectrum can not plausibly be produced by Comptonization of radiation from an accretion disk around the putative supermassive black hole in the center of Mrk 501: The spectrum emitted by an optically thick, geometrically thin accretion disk is reasonably well approximated by a blackbody spectrum whose temperature, for a black hole mass of $`10^8M_{}`$, yields an average photon energy of the disk radiation of $`ϵ_Dh\nu _D/(m_ec^2)10^5`$. If external Comptonization is to be efficient in competition with the synchrotron self-Compton process, the blob has to be rather close to the accretion disk, or a significant fraction of the disk photons has to be rescattered into the jet by surrounding material, so that the bulk of disk photons enters the blob from the side and is blue shifted by a factor of $`\mathrm{\Gamma }`$ into the blob rest frame. Thus, due to the strong reduction of the Klein-Nishina cross section for $`\gamma _eϵ^{}1`$ (where $`ϵ^{}`$ is the photon energy in the comoving frame), no significant radiative output at (observer’s frame) energies $`ϵ_{obs}D/(ϵ_D\mathrm{\Gamma })10^5`$ (where $`D`$ is the Doppler factor), corresponding to $`E_{obs}100`$ GeV, will be produced by external Comptonization, ruling out this process to explain the observed high-energy spectrum extending to TeV energies. If the high-energy spectrum in the $`100`$ GeV regime is produced by the SSC process, then the level of SSC radiation at 1 GeV may be estimated by
$$\nu F_\nu ^{SSC}(1\mathrm{GeV})\nu F_\nu ^{SSC}(E_{pk})\left(\frac{E_{pk}}{1\mathrm{GeV}}\right)^{\frac{p3}{2}}$$
(16)
where $`E_{pk}0.1`$ TeV is the energy of the $`\nu F_\nu `$ peak in the high-energy part of the spectrum and $`p`$ is the spectral index of the injected electron spectrum. Inserting a typical value of $`\nu F_\nu (E_{pk})10^{13}\mathrm{Jy}\mathrm{Hz}`$ and $`p2.5`$, this yields $`\nu F_\nu ^{SSC}(1\mathrm{GeV})310^{12}\mathrm{Jy}\mathrm{Hz}`$, which is already close to the maximum ever observed of $`10^{13}\mathrm{Jy}\mathrm{Hz}`$. In reality, the SSC spectrum is not a straight power-law below $`E_{pk}`$, but shows a gradual turnover so that the actual SSC flux at 1 GeV is substantially higher than the above estimate, leaving little room for an additional EIC contribution. For our fitting procedure, we are thus using a simplified version of our jet simulation code, in which the photon output (but not the electron cooling) from Comptonization of accretion disk radiation is neglected. These simulations properly account for the self-consistent cooling of the electron population and $`\gamma \gamma `$ absorption and pair production in the blob. In the simulations, the exact, angle-averaged synchrotron spectrum of an isotropic electron population as given by Crusius & Schlickeiser (crusius86 (1986)) is used.
Fig. 4 shows two examples of fits to low and high flux levels of Mkn 501. The respective fit parameters are given in the figure captions. Inverse-Compton cooling on accretion-disk photons may be important close to the base of the jet even if the resulting photon spectra do not contribute significantly to the time-averaged emission and is thus still included in our simulations.
Due to the non-linear nature of the model system, the fit results cannot be described in a simple, linear way as a function of the model parameters. We therefore construct a three-dimensional mesh of simulations in parameter space, with the electron density $`n_e`$, the high-energy cutoff $`\gamma _2`$ and the spectral index $`p`$ of the injected electron spectrum as parameters which are free to vary on the grid points.
We calculate our parameter grid varying $`p`$ in steps of 0.025 between 1.6 and 2.8, the total electron number density $`n_e`$ in steps of 5 cm<sup>-3</sup> between 10 and 120 cm<sup>-3</sup>, and $`\gamma _2`$ for values of $`210^6`$, $`310^6`$, $`510^6`$, $`7.510^6`$, $`10^7`$, $`1.510^7`$, $`210^7`$, $`2.510^7`$, and $`310^7`$. We constrain the range of $`\gamma _2`$ values to $`\gamma _2310^7`$ because of the kinematic limit and because around this energy, electron cooling due to triplet pair production on the highest-energy synchrotron photons, which is ignored in our simulations, becomes dominant over Compton scattering (Mastichiadis et al. mastichiadis94 (1994), Anguelov et al. ang99 (1999)). We find that our simulated spectra are only very weakly dependent on the actual value of $`\gamma _2`$. A change of $`\gamma _2`$ by a factor 3 typically results in an increase of the reduced $`\chi ^2`$ of only 0.3 such that the above-mentioned restriction of $`\gamma _2`$ values has only minor impact on our results. In fact $`\gamma _2`$ can be regarded as constant and of the order of $`10^7`$. We point out that the instantaneous synchrotron spectra of individual blobs at the time of injection, which might correspond to short-term X-ray and TeV flares, have their synchrotron peak at $`\nu _{sy,inst.}2.810^6(B/\mathrm{G})D\gamma _2^2`$ Hz if $`p<3`$, in agreement with the shift of the synchrotron peak into the hard X-ray regime during extreme flaring activity (e. g., Pian et al. pian (1998)).
All other parameters (in particular, the magnetic field at the particle injection site, $`B_0=0.05`$ G, the low-energy cutoff of the electron spectrum $`\gamma _1=300`$, and the Doppler factor, $`D=30`$) are fixed to values allowing good fits to the observed weekly averaged SEDs using our simulation code. An estimate for the required parameters can be found on the basis of the location of the synchrotron and SSC peaks of the observed broadband spectra as described below.
Although we are assuming the injection of a single power-law distribution of ultrarelativistic electrons into the jet, the time-averaged radiation spectrum will be reasonably well approximated by the one produced by a broken power-law distribution of electrons with spectral index $`p`$ below the break energy $`\gamma _b`$, and $`p+1`$ above the break energy. $`\gamma _b`$ may be computed by setting the synchrotron cooling time scale equal to the dynamical time scale of jet evolution (magnetic field decay), which yields
$$\gamma _b6.410^5\frac{\mathrm{\Gamma }_{25}}{z_{0.03}B_1^2}$$
(17)
where $`\mathrm{\Gamma }_{25}\mathrm{\Gamma }/25`$, the height of the injection/acceleration site above the accretion disk is $`z_i=0.03z_{0.03}`$ pc, and $`B_1=B_0/(0.1\mathrm{G})`$. Thus, our model calculations will produce a time-averaged synchrotron break at
$$\nu _{sy}3.410^{18}\overline{B}_1D_{30}\frac{\mathrm{\Gamma }_{25}^2}{z_{0.03}^2B_1^4}\mathrm{Hz}$$
(18)
where $`\overline{B}_1`$ is an appropriate average of the magnetic field (in units of $`0.1`$ G) over the jet evolution. For the purpose of these estimates, we neglect factors of $`(1+z)1`$ for Mkn 501. The location of the peak of the SSC component will be strongly influenced by Klein-Nishina effects and will thus depend on the actual shape of the synchrotron spectrum, which, in turn, depends on the electron spectral index $`p`$. Considering these effects, Tavecchio et al. (tav98 (1998)) find
$$ϵ_{SSC}\gamma _bDg(\alpha _1,\alpha _2)$$
(19)
where $`\alpha _1=(p1)/2`$, $`\alpha _2=p/2`$, and
$$g(\alpha _1,\alpha _2)=\mathrm{exp}\left[\frac{1}{\alpha _11}+\frac{1}{2(\alpha _2\alpha _1)}\right].$$
(20)
For $`p=2.5`$, this yields $`ϵ_{SSC}9.410^5D_{30}\mathrm{\Gamma }_{25}/(z_{0.03}B_1^2)`$, corresponding to
$$E_{SSC}490\frac{D_{30}\mathrm{\Gamma }_{25}}{z_{0.03}B_1^2}\mathrm{GeV}.$$
(21)
Combining Eqs. 18 and 21 and using the average observed $`ϵ_{sy}10^2`$ and $`ϵ_{SSC}10^6`$, we find
$$\frac{\overline{B}_1}{D_{30}}300\frac{ϵ_{sy}}{ϵ_{ssc}^2}g^2(\alpha _1,\alpha _2)0.34$$
(22)
for $`p=2.5`$ (see also Tavecchio et al. tav98 (1998)). Similarly, we may use Eq. (22) of Tavecchio et al. (tav98 (1998)) to estimate
$$\overline{B}D^{2+\alpha _1}\left[g(\alpha _1,\alpha _2)ϵ_{SSC}ϵ_{sy}\right]^{(1\alpha _1)/2}\sqrt{\frac{2f(\alpha _1,\alpha _2)}{c^3}}\frac{\left(\nu L_\nu \right)_{sy}}{t_{var}\sqrt{\left(\nu L_\nu \right)_{SSC}}}$$
(23)
where $`f(\alpha _1,\alpha _2)=1/(1\alpha _1)+1/(\alpha _21)`$. Using $`p=2.5`$, $`\left(\nu L_\nu \right)_{sy}610^{43}`$ erg/s, $`\left(\nu L_\nu \right)_{SSC}210^{44}`$ erg/s, and $`t_{var}5`$ h (Aharonian et al. 1999a ), we find
$$\overline{B}_1D_{30}^{11/4}0.34.$$
(24)
Combining this with Eq. 22, we have $`D_{30}1`$ and $`\overline{B}_10.34`$. These numbers are consistent with the limits found by Bednarek & Protheroe bednarek99 (1999), but are slightly outside the allowed region of parameter space as found by Tavecchio et al. (tav98 (1998)) and Kataoka et al. (kataoka (1999)). This is because in those papers the broadband spectrum is either characterized by quantities pertaining to individual outbursts or to a long-term quiescent state. Those parameters are not representative of the weekly averages investigated in this paper.
Having constructed the three-dimensional mesh of simulations, we compare all weekly SEDs with the simulated spectra and find the simulation with the smallest $`\chi ^2`$. Results of this procedure are described in the next section.
## 3 Results
### 3.1 Correlation TeV-X-Ray
The SSC model for Mkn 501 predicts a very strong correlation between the emission at the synchrotron peak (in soft – hard X-rays) and at the inverse-Compton peak (close to TeV $`\gamma `$-ray energies). We have derived an analytic estimate for the expected correlation, for variations of several input parameters. In the following discussion, unprimed quantities are measured in the co-moving frame, while a superscript $``$ refers to quantities measured in the observer’s frame. We assume that the time-averaged (cooling) electron spectrum can be described by a broken power-law,
$$n_e(\gamma )=n_0\{\begin{array}{cc}(\gamma /\gamma _b)^p\hfill & \text{for }\gamma _1\gamma \gamma _b\hfill \\ (\gamma /\gamma _b)^{(1+p)}\hfill & \text{for }\gamma _b\gamma \gamma _2\hfill \end{array}$$
(25)
where the injection spectral index $`2<p<3`$, and $`\gamma _b`$ is the break energy of the spectrum, determined by Eq. 17. The normalization is given by $`n_0n_e\gamma _b^p\gamma _1^{p1}/(p1)`$. We are using a $`\delta `$ approximation for the synchrotron spectrum:
$$L_{sy}(ϵ)=L_0\{\begin{array}{cc}(ϵ/ϵ_b)^{\frac{1p}{2}}\hfill & \text{for }ϵ_1ϵϵ_b\hfill \\ (ϵ/ϵ_b)^{\frac{p}{2}}\hfill & \text{for }ϵ_bϵϵ_2\text{,}\hfill \end{array}$$
(26)
where $`ϵ=h\nu /(m_ec^2)`$ is the dimensionless photon energy and $`ϵ_i=2.310^{14}(B/\mathrm{G})\gamma _i^2`$ is the characteristic synchrotron energy radiated by an electron of Lorentz factor $`\gamma _i`$. Normalizing the synchrotron luminosity to
$$L_{sy}B^2\underset{\gamma _1}{\overset{\gamma _2}{}}𝑑\gamma n_e(\gamma )\gamma ^2,$$
(27)
we have
$$L_0\frac{Bn_e}{p1}\left(\frac{\gamma _1}{\gamma _b}\right)^{p1}.$$
(28)
Neglecting Compton scattering events in the Klein-Nishina regime, $`ϵ\gamma >3/4`$, we may approximate the SSC spectrum by
$$L_{SSC}(ϵ_s)\underset{ϵ_1}{\overset{ϵ_2}{}}𝑑ϵ\frac{L_{sy}(ϵ)}{ϵ}\sqrt{\frac{ϵ_s}{ϵ}}n_e\left(\sqrt{\frac{3ϵ_s}{4ϵ}}\right)\mathrm{\Theta }\left(3/4\sqrt{ϵ_sϵ}\right),$$
(29)
where $`\mathrm{\Theta }`$ is the Heaviside function. The evaluation of this expressions is straightforward. Observed fluxes in the ASM, BATSE, and HEGRA energy ranges are calculated integrating the Doppler boosted synchrotron and SSC spectra, $`L^{}(ϵ^{})=D^3L(ϵ^{}/D)`$, over the energy ranges $`410^3ϵ_{ASM}^{}210^2`$, $`410^2ϵ_{BATSE}^{}0.4`$, and $`310^6ϵ_{HEGRA}^{}610^7`$.
In Fig. 5, we plot trajectories in the $`(F_{BATSE},F_{HEGRA})`$ and $`(F_{ASM},F_{HEGRA})`$ planes of these solutions, varying individual model parameter separately while all others are held constant at values representative of states of moderate X-ray and high-energy $`\gamma `$-ray fluxes.
A variation of the electron density obviously yields a relation $`F_{SSC}F_{sy}^2`$ since the synchrotron flux depends linearly, the SSC flux quadratically on $`n_e`$. Note that this dependence may be altered due to an increasing $`\gamma \gamma `$ absorption opacity intrinsic to the source, which is not included in the analytical estimate (29) used to compute the HEGRA flux. A variation of the electron injection spectral index $`p`$ results in a relation which may be approximated by $`F_{HEGRA}F_{BATSE}^{1.4}`$ and $`F_{HEGRA}F_{ASM}^{1.6}`$, i. e. the dependence is weaker than quadratic.
A variation of the magnetic field strength leads to more complex flux variations due to the back-reaction on the break Lorentz factor $`\gamma _b`$ as a result of radiative cooling. For relatively strong magnetic fields ($`B0.3`$ G), the X-ray and TeV $`\gamma `$-ray fluxes become anti-correlated.
Finally, if the variability of this source were dominated by a variation of the bulk Lorentz factor $`\mathrm{\Gamma }`$, the X-ray and high-energy $`\gamma `$-ray fluxes would be expected to be approximately linearly correlated (the back-reaction on $`\gamma _b`$ leads to a slight deviation from a strictly linear correlation), as long as the observer is located within the $`1/\mathrm{\Gamma }`$ beaming cone of the jet. If $`\mathrm{\Gamma }`$ increases beyond $`1/\theta _{obs}`$, both the X-ray and TeV $`\gamma `$-ray fluxes start to decrease with increasing $`\mathrm{\Gamma }`$. The same quasi-linear correlation would be expected if the variability were due to a bending jet, i. e. a variation of $`\theta _{obs}`$.
The empirical correlation of the TeV and the X-ray emission of Mkn 501 in 1997 has already been studied extensively using the CT1/CT2 and the CTS data from HEGRA and the soft X-ray data from RXTE ASM: Aharonian et al. (1999b ) find the correlation coefficient for the daily averages to be $`0.61\pm 0.057`$. This maximum correlation is found for zero time-lag. We examine the correlation between the weekly RXTE, BATSE and HEGRA points from table 2. Fig. 6 shows the observed correlation between the HEGRA and BATSE measurements, fitted with a second-order polynomial as well as with a power-law with index 1.4, which is the theoretical prediction if the variability is caused solely by variations of the electron injection spectral index $`p`$. Both fits give acceptable values for the reduced $`\chi ^2`$ (1.1 and 1.65 respectively). Still, due to the large error bars, we can not confidently distinguish between these and similar correlations on the basis of the currently available data. In any case, there is no indication for a super-quadratic dependence between the X-ray and TeV fluxes, which would be inconsistent with a pure SSC model, unless there is a persistent quiescent level of emission above which the observed flaring behaviour is superimposed.
Figure 7 shows the correlation between the weekly TeV and Soft-X-ray points. The linear correlation coefficient is $`0.59`$, nearly the same as found by Aharonian et al. (1999b ) who compare the same data on a daily basis. The constant term of the linear fit is still consistent with zero. However, the reduced $`\chi ^2`$ of 4.7 is too large for a good fit. This is also the case for a fitted power-law with index 1.6 which gives $`\chi ^2=4.9`$. The systematic differences in time-coverage which are not taken into account in the determination of the error bars may be responsible for this. Still, the large linear correlation coefficient suggests that the correlation is nearly linear.
Figure 8 shows the correlation between the weekly Hard-X-ray and Soft-X-ray points. In this case, the time coverage is the same for both instruments, only the duty cycles are different. The correlation coefficient is 0.53 corresponding to a 0.5 % chance probability for a linear correlation. The constant term of the linear fit is very well consistent with zero. This figure also illustrates the difference in the dynamical ranges of the variability in the soft and the hard X-ray band. At BATSE energies, which are believed to be near the high-energy end of the synchrotron spectrum, the variability amplitude is about 50 % larger than at RXTE energies. This fits nicely into the scheme that the strongest variability takes place at the high-energy ends of both spectral components.
### 3.2 Model fit results
Each weekly SED is compared to a three-dimensional mesh of $`48\times 22\times 9=9504`$ simulations (48 different values of $`p`$, 22 values of $`n_e`$, and 9 values of $`\gamma _2`$), selecting the simulated broadband spectrum with the smallest $`\chi ^2`$. The resulting best-fit parameters are listed in table 3, along with the resulting $`\chi ^2`$ divided by the number of data points. Since we do not have a continuous sequence in parameter space, the actual number of degrees of freedom is questionable so that we use the above quantity to assess the quality of the fit.
Our best-fit parameters for each individual weekly averaged broadband spectrum are listed in Tab. 3. With few exceptions (MJD 50549.0, 50576.9, 50618.8, 50696.1), all fits resulted in acceptable $`\chi ^2`$ values.
We find a correlation between the electron injection spectral index $`p`$ and the X-ray and high-energy $`\gamma `$-ray fluxes, while there is no obvious correlation with the total electron density and/or the high-energy cut-off $`\gamma _2`$ of the injected electron spectrum. There also appears to be a weak anti-correlation between the jet filling factor $`f(\mathrm{\Delta }t_{rep})^1`$ and the X-ray and HEGRA fluxes. This could indicate that during states of relatively low activity, the fluxes are dominated by a quasi-steady component from a continuous jet, while in high-activity state the emission is dominated by more isolated, eruptive events. However, this latter correlation is much less pronounced than the correlation with the electron spectral index and will need to be tested on the basis of future, more sensitive observations.
In Fig. 9 the temporal variation of the best-fit values of $`p`$ and $`\mathrm{\Delta }t_{rep}`$ are compared to the variations of the soft and hard X-ray fluxes and the 1.5 TeV flux. The correlation between the hard X-ray and TeV $`\gamma `$-ray fluxes with the injection spectral index is illustrated in figure 10.
Furthermore, we show in figure 11 the correlation between $`p`$ and the positions of the peaks in the synchrotron and the inverse Compton component of the SED. The peak positions are not fit parameters but were determined by finding the local maxima in the weekly SEDs. We find that, in our model, there is a strong correlation between the peak positions and $`p`$ such that these parameters can be regarded as nearly identical. However, while the peak positions are directly observable, the electron spectral index is the more fundamental quantity.
Our results indicate that medium-timescale high activity states in X-rays and high-energy $`\gamma `$-rays are consistent with a hardening of the electron spectrum injected at the base of the jet. As pointed out in the previous subsection, a pure SSC model in which the TeV $`\gamma `$-ray flux is strongly influenced by Klein-Nishina effects, predicts that the HEGRA and both the soft and hard X-ray fluxes should roughly be correlated by power-laws $`F_{HEGRA}F_X^\delta `$ with $`1.4\delta 1.6`$ in high flux-level states, in which the contribution of a possible quasi-stationary radiation component is small. The data available for this study are consistent with this but do not allow a clear distinction between different variability mechanisms, and future observations with increased sensitivity, in particular at multi-GeV to TeV energies, are needed in order to test this prediction.
## 4 Summary and Conclusions
We have presented broadband spectra of the extreme HBL Mkn 501 during its high state in 1997, including radio, optical, soft and hard X-ray, and TeV $`\gamma `$-ray observations. In this study we concentrated on the medium-timescale variability, using weekly averaged SEDs. We confirmed the strong correlation between the TeV $`\gamma `$-ray flux and the hard X-ray flux. This correlation was found to be non-linear and could be fitted with a second-order polynomial, in agreement with the expectation of an SSC dominated leptonic jet model, if the flux variations are related to fluctuations of the electron density in the jet and/or the spectral index of the electron spectrum at the time of injection into the jet.
The weekly averaged SEDs were fitted with a leptonic jet model, strongly dominated by the SSC process. With a few exceptions, this model yielded acceptable fits to the observed broadband spectra. The observed spectral variability of Mkn 501 could be explained mainly by variations of the electron spectral index. No clear correlation between the maximum electron energy and the hard X-ray and TeV $`\gamma `$-ray fluxes on the 1-week timescale was found, in contrast to the short-term variability of Mkn 501. Pian et al. (pian (1998)) have shown that the intraday variability of this object is most probably related to an increase of $`\gamma _2`$, leading to pronounced flares in hard X-rays, most probably on the synchrotron cooling timescale which is most likely $``$ a few hours and thus much shorter than the 1-week timescale considered in this paper. Our result indicates that such synchrotron flares are isolated events and are at most weakly correlated to the activity of the source on the 1-week timescale.
Our result that the flaring behaviour on intermediate timescales is consistent with a hardening of the electron spectrum is in contrast to the flaring characteristics observed in quasars. Recently, Mukherjee et al. (mukherjee99 (1999)) have investigated all available broadband data on the very luminous flat-spectrum radio quasar (FSRQ) PKS 0528+134, and found that its flaring behaviour is consistent with an increasing contribution of the external inverse-Compton component during flares, possibly related to an increase in the bulk Lorentz factor. The fits to the SEDs of PKS 0528+134 required that the average energy of relativistic electrons in the jet shifts towards lower values during flares, in contrast to the results found for Mkn 501. As pointed out by Böttcher (boettcher99 (1999)), this implies that the synchrotron peak is expected to shift towards lower frequencies during flares of FSRQs, while Mkn 501 and Mkn 421 show clear evidence for a shift of the synchrotron peak to higher frequencies.
We point out that in the present study the magnetic field along the jet and the bulk Lorentz factor of individual blobs were fixed, so that we cannot confidently rule out variations of the Doppler factor, accompanied by appropriate changes of the electron injection spectrum, as the flaring mechanism for Mkn 501. However, the very moderate variability at optical frequencies, as observed in Mkn 501, leads us to consider this flaring mechanism less likely in this object since it would require a peculiar conspiracy between the Doppler factor and the electron spectrum to keep the optical flux at an approximately constant level.
## Acknowledgements
The work of DP is supported by the Spanish CICYT grant SB97-B12601316. The work of MB was supported by NASA grant NAG 5-4055 (until Aug. 1999) and by Chandra Postdoctoral Fellowship grant number PF 9-10007, awarded by the Chandra X-ray Center, which is operated by the Smithsonian Astrophysical Observatory for NASA under contract NAS 8-39073. The RXTE ASM data has been obtained through the High Energy Astrophysics Science Archive Research Center Online Service provided by the NASA/Goddard Space Flight Center. We thank C.D. Dermer for valuable comments on the manuscript and J.M. Holeczek for providing a fit routine. |
warning/0002/math0002165.html | ar5iv | text | # Degenerate affine Hecke algebras and centralizer construction for the symmetric groups
## 1 Introduction
The centralizer construction proposed in shows that certain “quantum” algebras can be obtained as projective limits of centralizers in classical enveloping algebras. This approach has been applied to the series of matrix Lie algebras to construct the quantum algebras called the Yangians and twisted Yangians which are originally defined as certain deformations of enveloping algebras; cf. . The type $`A`$ case is treated in , and extended to the $`B,C,D`$ types in and . A modified version of the $`A`$ type construction is given in .
In the $`A`$ type case, one considers the centralizer $`𝒜_m(n)`$ of the subalgebra $`𝔤𝔩(nm)`$ in the enveloping algebra $`\mathrm{U}(𝔤𝔩(n))`$. It turns out that for each pair $`n>m`$ there is a natural algebra homomorphism $`𝒜_m(n)𝒜_m(n1)`$ so that one can define a projective limit algebra $`𝒜_m`$ by using the chain of homomorphisms
$$\mathrm{}𝒜_m(n)𝒜_m(n1)\mathrm{}𝒜_m(m+1)𝒜_m(m).$$
(1.1)
The algebra $`𝒜_m`$ has a large center $`𝒜_0`$ which is isomorphic to the algebra of shifted symmetric functions $`\mathrm{\Lambda }^{}`$ (see ) and one has an isomorphism
$$𝒜_m𝒜_0\mathrm{Y}_m,$$
(1.2)
where $`\mathrm{Y}_m`$ is the Yangian for the Lie algebra $`𝔤𝔩(m)`$; \[30, Theorem 2.1.15\]. In particular, for each $`nm`$ there is a natural homomorphism
$$\mathrm{Y}_m𝒜_m(n).$$
(1.3)
This result was used in to study the class of representations of $`\mathrm{Y}_m`$ which naturally arises from this construction. A similar result for the $`B,C,D`$ types was used in to construct weight bases of Gelfand–Tsetlin type for representations of the classical Lie algebras.
Our aim in the present paper is to extend these constructions to the series of the symmetric groups $`S(n)`$. Denote by $`B_m(n)`$ the centralizer of the subgroup $`S(nm)`$ in the group algebra $`[S(n)]`$, where $`S(nm)`$ consists of the permutations which fix each of the indices $`1,2,\mathrm{},m`$. It turns out that no natural analog of the chain (1.1) exists for the algebras $`B_m(n)`$. Indeed, note that $`B_n(n)=B_{n1}(n)=[S(n)]`$ and so, by analogy with (1.1) we would have a homomorphism $`[S(n)][S(n1)]`$ identical on $`[S(n1)]`$. But such a homomorphism does not exist for $`n>4`$.
On the other hand, it was observed in that an analog of (1.3) still exists: for any $`nm`$ there is a homomorphism $`_mB_m(n)`$, where $`_m`$ is the degenerate affine Hecke algebra introduced by Drinfeld and Lusztig . This fact was used by Okounkov and Vershik to develop a new approach to the representation theory of the symmetric groups; see also earlier results by Cherednik .
This observation together with the semigroup method allows one to expect that an analog of the centralizer construction for the symmetric group should exist, with the algebras $`[S(n)]`$ replaced with the semigroup algebras $`A(n)=[\mathrm{\Gamma }(n)]`$, where $`\mathrm{\Gamma }(n)`$ is the semigroup of partial bijections of the set $`\{1,\mathrm{},n\}`$. Alternatively, the elements of $`\mathrm{\Gamma }(n)`$ can be identified with $`(0,1)`$-matrices which have at most one $`1`$ in each row and column. The semigroups $`\mathrm{\Gamma }(n)`$ are studied in and used to prove Lieberman’s classification theorem for unitary representations of the complete infinite symmetric group.
Taking the centralizer $`A_m(n)`$ of $`\mathrm{\Gamma }(nm)`$ in $`A(n)`$ instead of the algebras $`B_m(n)`$ we do obtain an analog of the chain of homomorphisms (1.1) for the algebras $`A_m(n)`$. The corresponding projective limit algebra $`A_m`$ has a decomposition analogous to (1.2)
$$A_mA_0\stackrel{~}{}_m,$$
(1.4)
where $`\stackrel{~}{}_m`$ is a ‘semigroup analog’ of the degenerate affine Hecke algebra $`_m`$. The algebra $`\stackrel{~}{}_m`$ can be presented by generators and defining relations. Moreover, the algebra $`_m`$ is a homomorphic image of $`\stackrel{~}{}_m`$. The subalgebra $`A_0A_m`$ is commutative and isomorphic to the algebra of shifted symmetric functions $`\mathrm{\Lambda }^{}`$; see for a detailed study of the algebra $`\mathrm{\Lambda }^{}`$.
The mentioned above homomorphisms $`_mB_m(n)`$ can be regarded as ‘retractions’ of the homomorphisms $`\stackrel{~}{}_mA_m(n)`$ whose existence is provided with the centralizer construction for the symmetric groups.
Finally, the algebra $`A`$ is defined as the inductive limit of $`A_m`$ as $`m\mathrm{}`$. We show that $`A`$ naturally acts in the so-called tame representations of $`S(\mathrm{})`$ and it can be regarded as the ‘true’ analog of the group algebra $`[S(\mathrm{})]`$. Indeed, contrary to the finite case of $`[S(n)]`$, the algebra $`[S(\mathrm{})]`$ has a trivial center while $`A`$ contains a large center whose elements act by scalar operators in the tame representations of $`S(\mathrm{})`$. Moreover, the central elements separate the irreducible tame representations.
Some other generalizations of the degenerate affine Hecke algebra $`_m`$ have been constructed by Nazarov .
The paper is organized as follows. Section 2 is preliminary. Here we define tame representations of $`S(\mathrm{})`$ and describe the properties of the semigroups $`\mathrm{\Gamma }(n)`$; most of these results are contained in . In Section 3 we construct the algebras $`A_m`$ as projective limits of the centralizers $`A_m(n)`$. Section 4 describes the algebra $`A_0`$ and establishes its isomorphism with the algebra of shifted symmetric functions $`\mathrm{\Lambda }^{}`$. The main results are given in Section 5 where we investigate the structure of $`A_m`$ and describe its relationship with the degenerate affine Hecke algebras.
The authors acknowledge the financial support of the Australian Research Council. G.O. was also supported by the Russian Foundation for Basic Research under grant 98-01-00303.
## 2 Tame representations and semigroups
Here we give constructions of the tame representations of the infinite symmetric group and describe the semigroups of partial bijections. The material of this section is based on and \[27, §2\]. About applications of the semigroup method to representations of “big” groups, see Olshanski , Neretin .
### 2.1 Constructing tame representations
Let $`=\{1,2,\mathrm{}\}`$ and $`n`$. We denote by $`S(n)`$ the group of permutations of the set $`_n:=\{1,\mathrm{},n\}`$. We also regard $`S(n)`$ as the group of permutations of $``$ fixing $`n+1,n+2,\mathrm{}`$, and set
$$S(\mathrm{})=\underset{n1}{}S(n).$$
This is the group of all finite permutations of the set $``$.
For any $`mn`$, we denote by $`S_m(n)`$ the subgroup of permutations in $`S(n)`$ fixing $`1,\mathrm{},m`$, and set
$$S_m(\mathrm{})=\underset{nm}{}S_m(n).$$
Note that the subgroups $`S(n)`$ and $`S_n(\mathrm{})`$ of $`S(\mathrm{})`$ commute. This simple observation will play an important role later.
For a (unitary) representation $`T`$, we denote by $`H(T)`$ its (Hilbert) space.
Let $`T`$ be a unitary representation of the group $`S(\mathrm{})`$. For $`n=0,1,\mathrm{}`$, denote by $`H_n(T)`$ the subspace of $`S_n(\mathrm{})`$-invariant vectors in $`H(T)`$. Since $`\{S_n(\mathrm{})\}`$ is an descending chain of subgroups, $`\{H_n(T)\}`$ is an ascending chain of subspaces. Let
$$H_{\mathrm{}}(T)=\underset{n0}{}H_n(T).$$
Since $`S(n)`$ and $`S_n(\mathrm{})`$ commute, the subspace $`H_n(T)`$ is invariant with respect to $`S(n)`$, and so, $`H_{\mathrm{}}(T)`$ is invariant with respect to the whole group $`S(\mathrm{})`$.
###### Definition 2.1
A unitary representation $`T`$ of the group $`S(\mathrm{})`$ is said to be tame if $`H_{\mathrm{}}(T)`$ is dense in $`H(T)`$.
Note that for an irreducible representation $`T`$, this is equivalent to saying that $`H_{\mathrm{}}(T)`$ is nonzero.
Clearly, the trivial representation of $`S(\mathrm{})`$ is tame, and another one–dimensional representation, $`s\mathrm{sgn}s`$, is not tame. Less trivial examples follow.
###### Example 2.2
(i) Let $`H`$ be the space $`l_2`$ with its canonical basis $`e_1,e_2,\mathrm{}`$ and let $`S(\mathrm{})`$ operate in $`H`$ by permuting the basis vectors. The representation $`H`$ is tame. It is irreducible and for any $`n`$ the subspace $`H_n`$ is spanned by $`e_1,\mathrm{},e_n`$.
(ii) The right (or left) regular representation of the group $`S(\mathrm{})`$ in the Hilbert space $`l_2\left(S(\mathrm{})\right)`$ is not tame.
For any $`n=0,1,2,\mathrm{}`$ and any partition $`\lambda n`$, we will construct a tame representation $`T_\lambda `$. First, if $`n=0`$ then $`\lambda =\mathrm{}`$ and $`T_{\mathrm{}}`$ is the one-dimensional trivial representation. Let now $`n1`$ and let $`\pi _\lambda `$ denote the irreducible representation of $`S(n)`$ corresponding to $`\lambda `$. Then $`T_\lambda `$ is defined as the induced representation
$$T_\lambda =\mathrm{Ind}_{S(n)\times S_n(\mathrm{})}^{S(\mathrm{})}(\pi _\lambda 1)$$
(2.1)
where 1 stands for the trivial representation of $`S_n(\mathrm{})`$. The representation $`T_\lambda `$ can be realized as follows. Let $`\mathrm{\Omega }(n)`$ denote the set of injective maps $`\omega :_n.`$ Define a right action of $`S(\mathrm{})`$ on $`\mathrm{\Omega }(n)`$ by
$$\omega s=s^1\omega ,sS(\mathrm{}),$$
(2.2)
and a left action of $`S(n)`$ by
$$t\omega =\omega t^1,tS(n).$$
(2.3)
Note that these two actions commute. Consider the Hilbert space $`l_2(\mathrm{\Omega }(n),H(\pi _\lambda ))`$ of $`H(\pi _\lambda )`$-valued square-integrable functions on $`\mathrm{\Omega }(n)`$, and let $`H(n,\lambda )`$ be its subspace formed by the functions $`f(\omega )`$ such that
$$f(t\omega )=\pi _\lambda (t)f(\omega ),tS(n),\omega \mathrm{\Omega }(n).$$
(2.4)
The action of $`T_\lambda `$ in $`H(n,\lambda )`$ is given by
$$(T_\lambda (s)f)(\omega )=f(\omega s).$$
(2.5)
The space $`H(T_\lambda )`$ may now be identified with $`H(n,\lambda )`$. For any $`ln`$ set
$$\mathrm{\Omega }(n,l)=\{\omega \mathrm{\Omega }(n)|\omega (_n)_l\}$$
(2.6)
and note that
$$\mathrm{\Omega }(n)=\underset{ln}{}\mathrm{\Omega }(n,l)$$
(2.7)
Also, set
$$H_l^{}(T_\lambda )=\{fH(n,\lambda )|\mathrm{supp}f\mathrm{\Omega }(n,l)\},$$
(2.8)
where $`\mathrm{supp}f=\{\omega \mathrm{\Omega }(n)|f(\omega )0\}`$.
###### Proposition 2.3
For any $`n`$ we have
$$H_l(T_\lambda )=\{\begin{array}{cc}\{0\}\hfill & \text{if }l<n,\hfill \\ H_l^{}(T_\lambda )\hfill & \text{if }ln.\hfill \end{array}$$
(2.9)
Proof. Since $`S_l(\mathrm{})`$ acts trivially on $`\mathrm{\Omega }(n,l)`$, we have $`H_l^{}(T_\lambda )H_l(T_\lambda )`$. Conversely, let $`fH_l(T_\lambda )`$. Then the function $`f(\omega )^2`$ is constant on any orbit of the subgroup $`S_l(\mathrm{})`$ in $`\mathrm{\Omega }(n)`$. Since the sum of the $`f(\omega )^2`$ taken over $`\omega \mathrm{\Omega }(n)`$ must be finite, we have $`f(\omega )=0`$ unless the orbit containing $`\omega `$ is finite. But if $`\omega \mathrm{\Omega }(n,l)`$, then its orbit is clearly infinite, so that $`f(\omega )=0`$. This proves the opposite inclusion $`H_l(T_\lambda )H_l^{}(T_\lambda )`$.
###### Proposition 2.4
For any $`n`$ and any $`\lambda n`$, the representation $`T_\lambda `$ of $`S(\mathrm{})`$ is tame and irreducible.
Proof. By Proposition 2.3,
$$H_{\mathrm{}}(T_\lambda )=\underset{l}{}H_l(T_\lambda )=\underset{l}{}H_l^{}(T_\lambda ).$$
(2.10)
This is the subspace of finitely supported functions in $`H(T_\lambda )=H(n,\lambda )`$ which is clearly dense. So, $`T_\lambda `$ is tame.
The subspace $`H_n(T_\lambda )=H_n^{}(T_\lambda )`$ is both cyclic in $`H(T_\lambda )`$ and irreducible under the action of the subgroup $`S(n)`$. This implies that $`T_\lambda `$ is irreducible.
We shall identify any partition $`\lambda `$ with its Young diagram; see e.g. . We write $`|\lambda |=n`$ if $`\lambda `$ has $`n`$ boxes. Given two diagrams $`\lambda `$ and $`\mu `$ the notation $`\mu \lambda `$ will mean that $`\mu `$ can be obtained from $`\lambda `$ by removing one box, i.e. $`\mu \lambda `$ and $`|\mu |=|\lambda |1`$. An infinite tableau is defined as an infinite chain of diagrams
$$\tau =(\lambda ^{(1)}\lambda ^{(2)}\mathrm{}),|\lambda ^{(n)}|=n.$$
(2.11)
Two infinite tableaux will be called equivalent if the corresponding chains of diagrams differ in a finite number of places only.
Given an infinite tableau $`\tau `$, we may construct an inductive limit unitary representation $`\mathrm{\Pi }(\tau )`$ of the group $`S(\mathrm{})`$ as follows. By the branching rule for the symmetric groups (see e.g. , ), for any $`n`$ the representation $`\pi _{\lambda ^{(n)}}`$ occurs in the decomposition of $`\pi _{\lambda ^{(n+1)}}S(n)`$ with multiplicity one. Hence there is an infinite chain of embeddings
$$\pi _{\lambda ^{(1)}}\pi _{\lambda ^{(2)}}\mathrm{}$$
(2.12)
which are defined uniquely up to scalar multiples. So we may set
$$\mathrm{\Pi }(\tau )=\mathrm{lim}\mathrm{ind}\pi _{\lambda ^{(n)}},n\mathrm{}.$$
(2.13)
One can show that any $`\mathrm{\Pi }(\tau )`$ is irreducible (cf. \[26, Theorem 2.1\] and \[27, §2.7\]), and that $`\mathrm{\Pi }(\tau )`$ and $`\mathrm{\Pi }(\tau ^{})`$ are isomorphic if and only if $`\tau `$ and $`\tau ^{}`$ are equivalent.
This construction provides a large class of pairwise non-equivalent irreducible representations of the group $`S(\mathrm{})`$. We will be only interested in some special representations of this class. Let $`\lambda =(\lambda _1,\mathrm{},\lambda _l)`$ be an arbitrary diagram. Consider an infinite tableau $`\tau =(\lambda ^{(i)})`$ such that
$$\lambda ^{(i)}=(i|\lambda |,\lambda _1,\mathrm{},\lambda _l)\text{for}i|\lambda |+\lambda _1,$$
(2.14)
and set $`\mathrm{\Pi }_\lambda =\mathrm{\Pi }(\tau ).`$ The representation $`\mathrm{\Pi }_\lambda `$ is well defined since the equivalence class of $`\mathrm{\Pi }_\lambda `$ does not depend on the choice of $`\lambda ^{(i)}`$ for small values of $`i`$.
###### Proposition 2.5
The representations $`T_\lambda `$ and $`\mathrm{\Pi }_\lambda `$ are isomorphic for any $`\lambda `$.
Proof. For any $`ln+\lambda _1`$, the natural representation of $`S(l)`$ in the space $`H_l(T_\lambda )`$ is isomorphic to the induced representation
$$\mathrm{Ind}_{S(n)\times S(ln)}^{S(l)}(\pi _\lambda 1)$$
(2.15)
where $`S(ln)`$ is identified with $`S_n(l)`$, and $`1`$ stands for the trivial representation of $`S(ln)`$. This follows immediately from (2.9). It is well known that the representation (2.15) is multiplicity free and that its spectrum consists of the representations $`\pi _\mu `$ such that $`\mu l`$ and
$$\mu _{i+1}\lambda _i\mu _i,i1.$$
(2.16)
It follows from (2.16) that $`\pi _{\lambda ^{(l)}}`$ occurs in the decomposition of (2.15). Let $`H_l^0(T_\lambda )`$ denote the corresponding subspace of $`H_l(T_\lambda )`$. It remains to prove that $`H_m^0(T_\lambda )`$ is contained in $`H_l^0(T_\lambda )`$ for any $`m<l`$ provided that $`m`$ is large enough. This follows from the fact that
$$\pi _{\lambda ^{(m)}}\pi _{\lambda ^{(l)}}|_{S(m)},$$
(2.17)
and $`\pi _{\lambda ^{(m)}}\pi _\mu |_{S(m)}`$ if $`\mu `$ satisfies (2.16) and $`\mu \lambda ^{(l)}`$. Indeed, property (2.17) follows from the definition of the diagrams $`\lambda ^{(l)}`$ for large $`l`$ and the branching rule. Finally, note that there exists $`i`$ such that $`\lambda _i>\mu _{i+1}`$ (otherwise $`\mu =\lambda ^{(l)}`$). Applying the branching rule again we complete the proof.
### 2.2 The semigroup method
###### Definition 2.6
Let $`X`$ be a set.
(i) A partial bijection of $`X`$ is a bijection $`\gamma :DR`$ between two (possibly empty) subsets of $`X`$. The subset $`DX`$ is called the domain of $`\gamma `$ and denoted by $`\mathrm{dom}\gamma `$. The subset $`RX`$ is called the range of $`\gamma `$ and denoted by $`\mathrm{range}\gamma `$. If $`xX`$ belongs to $`\mathrm{dom}\gamma `$, then we will say that $`\gamma `$ is defined on $`x`$. The set of partial bijections of $`X`$ is denoted by $`\mathrm{\Gamma }(X)`$.
(ii) Given $`\gamma \mathrm{\Gamma }(X)`$, we define $`\gamma ^{}\mathrm{\Gamma }(X)`$ as the inverse bijection $`\gamma ^1:\mathrm{range}\gamma \mathrm{dom}\gamma `$, so that $`\mathrm{dom}\gamma ^{}=\mathrm{range}\gamma `$ and $`\mathrm{range}\gamma ^{}=\mathrm{dom}\gamma `$. The mapping $`\gamma \gamma ^{}`$ is involutive: $`(\gamma ^{})^{}=\gamma `$.
(iii) Given $`\gamma _1,\gamma _2\mathrm{\Gamma }(X)`$ with $`\mathrm{dom}\gamma _i=D_i`$ and $`\mathrm{range}\gamma _i=R_i`$, $`i=1,2`$, their product $`\gamma _1\gamma _2`$ is a partial bijection on $`X`$ with $`D=\mathrm{dom}\gamma _1\gamma _2=\gamma _2^1(D_1R_2)`$ and $`R=\mathrm{range}\gamma _1\gamma _2=\gamma _1(D_1R_2)`$:
$$\gamma _1\gamma _2=(\gamma _1|_{\gamma _2(D)})(\gamma _2|_D).$$
(2.18)
That is, $`\gamma _1\gamma _2`$ is defined on $`xX`$ if and only if $`\gamma _2`$ is defined on $`x`$ and $`\gamma _1`$ is defined on $`\gamma _2(x)`$; then $`(\gamma _1\gamma _2)(x)=\gamma _1(\gamma _2(x))`$.
Any $`\gamma \mathrm{\Gamma }(X)`$ may be regarded as a relation on $`X`$. The product defined above is the product of relations; see . With this product $`\mathrm{\Gamma }(X)`$ becomes a semigroup, called the semigroup of partial bijections of $`X`$.
The involution $`\gamma \gamma ^{}`$ is an anti-homomorphism of $`\mathrm{\Gamma }(X)`$, so that $`\mathrm{\Gamma }(X)`$ is an involutive semigroup.
The semigroup $`\mathrm{\Gamma }(X)`$ possesses the unity 1, which is the identity bijection $`XX`$, and the zero 0, which is the (unique) bijection of the empty subset of $`X`$ onto itself. Note that
$$1\gamma =\gamma \mathrm{\hspace{0.17em}1}=\gamma ,0\gamma =\gamma \mathrm{\hspace{0.17em}0}=0,\text{for any}\gamma \mathrm{\Gamma }(X).$$
(2.19)
For any subset $`YX`$, let $`1_Y\mathrm{\Gamma }(X)`$ denote the identity bijection $`YY`$. In particular, $`1_X=1`$ and $`1_{\mathrm{}}=0`$. Then $`1_Y`$ is a self-adjoint idempotent, i.e.,
$$(1_Y)^{}=1_Y,(1_Y)^2=1_Y.$$
(2.20)
Moreover, all idempotents of this type are pairwise commuting. For any $`\gamma \mathrm{\Gamma }(X)`$ we obviously have
$$\gamma ^{}\gamma =1_{\mathrm{dom}\gamma },\gamma \gamma ^{}=1_{\mathrm{range}\gamma }.$$
(2.21)
The subset of those $`\gamma \mathrm{\Gamma }(X)`$ for which $`\mathrm{dom}\gamma =\mathrm{range}\gamma =X`$ is clearly identified with the group $`S(X)`$ of all permutations of the set $`X`$.
Remark. The semigroup $`\mathrm{\Gamma }(X)`$ is a model example of an inverse semigroup (each of its elements has an inverse); see . The class of inverse semigroups is very closed to that of groups, and the role of the semigroups of partial bijections $`\mathrm{\Gamma }(X)`$ is quite similar to that of the symmetric groups $`S(X)`$. In particular, any inverse semigroup is isomorphic to an involutive subsemigroup of some $`\mathrm{\Gamma }(X)`$: this is an analog of Cayley theorem; see .
There is a convenient realization of partial bijections by (0,1)-matrices, i.e., by the matrices whose entries are 0 or 1. A (0,1)-matrix is said to be monomial if any of its rows or columns contains at most one 1. Given a set $`X`$, we will consider (0,1)-matrices whose rows and columns are labeled by the points of $`X`$. Then to any $`\gamma \mathrm{\Gamma }(X)`$, we assign a monomial (0,1)-matrix $`[\gamma _{xy}]`$ as follows: for $`x,yX`$
$$\gamma _{xy}=\{\begin{array}{cc}1\hfill & \text{if }y\mathrm{dom}\gamma \text{ and }\gamma (y)=x,\hfill \\ 0\hfill & \text{otherwise.}\hfill \end{array}$$
(2.22)
In particular, the unity $`1\mathrm{\Gamma }(X)`$ and the zero $`0\mathrm{\Gamma }(X)`$ are represented by the unit and the zero matrices, respectively. We thus obtain an isomorphism between the semigroup $`\mathrm{\Gamma }(X)`$ and the semigroup of monomial matrices over $`X`$ with the usual matrix multiplication. Note that in this matrix realization, the involution $`\gamma \gamma ^{}`$ is represented by the matrix transposition.
###### Definition 2.7
For $`YX`$ we define three mappings as follows:
(i) The canonical projection $`\theta :\mathrm{\Gamma }(X)\mathrm{\Gamma }(Y)`$. Let $`\gamma \mathrm{\Gamma }(X)`$. Then $`\theta (\gamma )`$ is defined at $`yY`$ if $`\gamma `$ is defined at $`y`$ and $`\gamma (y)Y`$. The image of $`y`$ with respect to $`\theta (\gamma )`$ is $`\gamma (y)`$. In matrix terms: the matrix of $`\theta (\gamma )`$ is obtained from that of $`\gamma `$ by striking the rows and columns corresponding to points of $`XY`$.
(ii) The canonical embedding $`\varphi :\mathrm{\Gamma }(Y)\mathrm{\Gamma }(X)`$. Let $`\gamma \mathrm{\Gamma }(Y)`$. Then $`\varphi (\gamma )`$ is defined at $`xX`$ if and only if $`x\mathrm{dom}\gamma Y`$ or $`xXY`$. In the former case $`\varphi (\gamma )`$ sends $`x`$ to $`\gamma (x)`$, and in the latter case it fixes $`x`$. In matrix terms, for $`x,yX`$,
$$\varphi (\gamma )_{xy}=\{\begin{array}{cc}\gamma _{xy}\hfill & \text{if }x,yY,\hfill \\ \delta _{xy}\hfill & \text{otherwise.}\hfill \end{array}$$
(2.23)
(iii) The $`0`$-embedding $`\psi :\mathrm{\Gamma }(Y)\mathrm{\Gamma }(X)`$. Let $`\gamma \mathrm{\Gamma }(Y)`$. Then $`\psi (\gamma )`$ is obtained by regarding $`\mathrm{dom}\gamma `$ and $`\mathrm{range}\gamma `$ as subsets of $`X`$. In matrix terms, for $`x,yX`$,
$$\psi (\gamma )_{xy}=\{\begin{array}{cc}\gamma _{xy}\hfill & \text{if }x,yY\text{,}\hfill \\ 0\hfill & \text{otherwise.}\hfill \end{array}$$
(2.24)
For a positive integer $`n`$ we shall denote by $`\mathrm{\Gamma }(n)`$ the semigroup $`\mathrm{\Gamma }(_n)`$. Given a tame representation $`T`$ of $`S(\mathrm{})`$ we now construct a representation $`𝒯_n`$ of the semigroup $`\mathrm{\Gamma }(n)`$. Regarding $`S(\mathrm{})`$ as the group of infinite monomial matrices which have only a finite number of $`1`$’s off the diagonal, define the map $`\theta ^{(n)}:S(\mathrm{})\mathrm{\Gamma }(n)`$ which assigns to each infinite matrix its upper left $`n\times n`$ submatrix. It is easy to check that the map $`\theta ^{(n)}`$ is surjective and it induces a bijection between the set of double cosets $`S_n(\mathrm{})\backslash S(\mathrm{})/S_n(\mathrm{})`$ and $`\mathrm{\Gamma }(n)`$; see .
Let $`T`$ be a tame representation of $`S(\mathrm{})`$. Denote by $`P_n`$ the orthogonal projection $`P_n:H(T)H_n(T)`$. Suppose that $`n`$ is so large that $`H_n(T)\{0\}`$. For any $`\sigma _1,\sigma _2S(\mathrm{})`$ we have
$$\theta ^{(n)}(\sigma _1)=\theta ^{(n)}(\sigma _2)P_nT(\sigma _1)P_n=P_nT(\sigma _2)P_n.$$
(2.25)
Therefore there exists a unique map $`𝒯_n:\mathrm{\Gamma }(n)\mathrm{End}H_n(T)`$ such that
$$𝒯_n\left(\theta ^{(n)}(\sigma )\right)=P_nT(\sigma )|_{H_n(T)},\sigma S(\mathrm{}).$$
(2.26)
###### Proposition 2.8
The map $`𝒯_n=𝒯_n(T)`$ defined by (2.26) is a representation of the semigroup $`\mathrm{\Gamma }(n)`$ in $`H_n(T)`$.
Proof. We give a sketch of the proof. The details and one more proof can be found in . We show first that the tame representation $`T`$ of $`S(\mathrm{})`$ can be extended to a representation $`𝒯`$ of the semigroup $`\mathrm{\Gamma }(\mathrm{})`$ of partial bijections of the set $``$.
Further, we prove that $`P_n`$ coincides with the operator $`𝒯(1_n)`$ where $`1_n`$ is the identity bijection of the subset $`_n`$.
Finally, consider the $`0`$-embedding $`\psi :\mathrm{\Gamma }(n)\mathrm{\Gamma }(\mathrm{})`$; see (2.24). We have for any $`\gamma \mathrm{\Gamma }(n)`$
$$𝒯_n(\gamma )=P_n𝒯\left(\psi (\gamma )\right)P_n=𝒯(1_n)𝒯\left(\psi (\gamma )\right)𝒯(1_n)=𝒯(1_n\psi (\gamma )1_n)=𝒯(\psi (\gamma )).$$
(2.27)
This proves the multiplicativity of $`𝒯_n`$.
Consider the tame representations $`T_\lambda `$ of $`S(\mathrm{})`$ constructed in the previous section. Proposition 2.8 yields a representation $`𝒯_n(\lambda ):=𝒯_n(T_\lambda )`$ of $`\mathrm{\Gamma }(n)`$ provided that $`H_n(T_\lambda )0`$, i.e., $`n|\lambda |`$.
We now outline the proof of the classification theorem for representations of $`\mathrm{\Gamma }(n)`$; see .
###### Theorem 2.9
The representations $`𝒯_n(\lambda )`$ where $`\lambda `$ is a partition with $`|\lambda |n`$ exhaust all irreducible representations of $`\mathrm{\Gamma }(n)`$.
Proof. Let $`𝒯`$ be a representation of $`\mathrm{\Gamma }(n)`$. Then for any $`mn`$ the operator $`𝒯(1_m)`$ is a projection. Denote its image by $`H_m(𝒯)`$. We let $`m(𝒯)`$ denote the minimum value of $`m`$ such that $`H_m(𝒯)\{0\}`$.
Further, if $`𝒯`$ is irreducible then one shows that $`dim𝒯<\mathrm{}`$. If $`m=m(𝒯)`$ then the subspace $`H_m(𝒯)`$ is invariant under $`S(m)`$ and irreducible. So, $`H_m(𝒯)`$ corresponds to a partition $`\lambda `$ with $`|\lambda |=m`$. A standard argument proves that $`𝒯`$ is uniquely determined by $`\lambda `$.
Conversely, given a partition $`\lambda `$ with $`|\lambda |=mn`$ one uses the following argument to construct an irreducible representation $`𝒯`$ of $`\mathrm{\Gamma }(n)`$ such that $`m=m(𝒯)`$ and the representation of $`S(m)`$ in $`H_m(𝒯)`$ corresponds to $`\lambda `$.
Denote by $`\mathrm{\Omega }(m,n)`$ the set of injective maps $`\omega :_m_n`$. We define $`H(𝒯)`$ to be the space of functions $`f:\mathrm{\Omega }(m,n)H(\pi _\lambda )`$ such that
$$f(t\omega )=\pi _\lambda (t)f(\omega ),tS(m),\omega \mathrm{\Omega }(m,n);$$
(2.28)
see (2.3). The action of $`\gamma \mathrm{\Gamma }(n)`$ is given by
$$(𝒯(\gamma )f)(\omega )=\{\begin{array}{cc}f(\gamma ^{}\omega )\hfill & \text{if }\omega =(\omega _1,\mathrm{},\omega _m)\mathrm{dom}\gamma ^{},\hfill \\ 0\hfill & \text{otherwise}.\hfill \end{array}$$
(2.29)
One easily checks that $`𝒯`$ is a representation of $`\mathrm{\Gamma }(n)`$. Moreover, it is isomorphic to the representation $`𝒯_n(\lambda )`$.
Note that the representation of $`\mathrm{\Gamma }(n)`$ corresponding to $`m=0`$ and the empty diagram is the trivial representation sending all elements of $`\mathrm{\Gamma }(n)`$ to $`1`$.
Theorem 2.9 leads to the following result.
###### Theorem 2.10
Let $`\lambda `$ range over the set of all Young diagrams including the empty diagram. The representations $`T_\lambda `$ constructed in Section 2.1 exaust, within equivalence, all the irreducible tame representations of the group $`S(\mathrm{})`$. Moreover, any tame representation of $`S(\mathrm{})`$ can be decomposed into a direct sum of irreducible tame representations.
Proof. See Theorem 6.7 and §7.2 in .
This is equivalent to Lieberman’s theorem concerning continuous unitary representations of the complete infinite symmetric group (this group consists of all permutations of the set $``$); see \[28, §7\].
## 3 Centralizer construction
We shall denote by $`\theta _n`$ the canonical projection $`\mathrm{\Gamma }(n)\mathrm{\Gamma }(n1)`$; see Definition 2.7. So, if $`\gamma \mathrm{\Gamma }(n)`$ then $`\theta _n(\gamma )`$ is the upper left corner of $`\gamma `$ of order $`n1`$. Here the elements of $`\mathrm{\Gamma }(n)`$ and $`\mathrm{\Gamma }(n1)`$ are regarded as $`(0,1)`$-matrices of order $`n`$ and $`n1`$, respectively.
We set $`A(n)=[\mathrm{\Gamma }(n)]`$, the semigroup algebra of $`\mathrm{\Gamma }(n)`$. The canonical embedding $`\mathrm{\Gamma }(n1)\mathrm{\Gamma }(n)`$ is extended to an embedding $`A(n1)A(n)`$ by linearity. Further, set $`A(\mathrm{})=[\mathrm{\Gamma }(\mathrm{})]=_{n1}A(n)`$, the semigroup algebra of $`\mathrm{\Gamma }(\mathrm{})`$.
For each $`i=1,\mathrm{},n`$ denote by $`\epsilon _i`$ the diagonal $`n\times n`$-matrix whose $`ii`$-th entry is $`0`$ and all other diagonal entries are equal to $`1`$. The corresponding element of $`\mathrm{\Gamma }(n)`$ is the identity bijection of the set $`_n\{i\}`$. The algebra $`A(n)`$ is obviously generated by $`S(n)`$ and the elements $`\epsilon _i`$, $`i=1,\mathrm{},n`$. We have for any $`i`$ and any $`\sigma S(n)`$:
$$\sigma \epsilon _i\sigma ^1=\epsilon _{\sigma (i)}.$$
(3.1)
###### Proposition 3.1
The algebra $`A(n)`$ is isomorphic to the abstract algebra with generators $`s_1,\mathrm{},s_{n1}`$, $`\epsilon _1,\mathrm{},\epsilon _n`$ and the defining relations
$`s_k^2`$ $`=1,`$ $`s_ks_{k+1}s_k`$ $`=s_{k+1}s_ks_{k+1},s_ks_l=s_ls_k,|kl|>1,`$ (3.2)
$`\epsilon _k^2`$ $`=\epsilon _k,`$ $`\epsilon _k\epsilon _l`$ $`=\epsilon _l\epsilon _k,`$ (3.3)
$`s_k\epsilon _k`$ $`=\epsilon _{k+1}s_k,`$ $`s_k\epsilon _k\epsilon _{k+1}`$ $`=\epsilon _k\epsilon _{k+1}.`$ (3.4)
Proof. Denote the abstract algebra by $`𝒜(n)`$. The assignments $`s_k(k,k+1)`$ and $`\epsilon _k\epsilon _k`$ obviously define an algebra epimorphism $`𝒜(n)A(n)`$. Note also that (3.2) are defining relations for the symmetric group $`S(n)`$ and (3.1) holds. To complete the proof we verify that $`dim𝒜(n)dimA(n)`$. We have
$$dimA(n)=|\mathrm{\Gamma }(n)|=\underset{r=0}{\overset{n}{}}\left(\genfrac{}{}{0pt}{}{n}{r}\right)^2r!.$$
(3.5)
On the other hand, we see from the relations (3.3) and (3.4) that
$$𝒜(n)=[S(n)][\epsilon _1,\mathrm{},\epsilon _n].$$
(3.6)
Here $`[\epsilon _1,\mathrm{},\epsilon _n]`$ is the subalgebra of $`𝒜(n)`$ generated by the $`\epsilon _k`$. It is spanned by the monomials $`\epsilon _{k_1}\mathrm{}\epsilon _{k_r}`$ with $`k_1<\mathrm{}<k_r`$. Given such a monomial, consider the subspace $`[S(n)]\epsilon _{k_1}\mathrm{}\epsilon _{k_r}`$ in $`𝒜(n)`$. Using (3.1) if necessary, we may assume without loss of generality that $`k_i=i`$ for each $`i`$. Observe that by (3.4) we have in $`𝒜(n)`$
$$\sigma S(r)\epsilon _1\mathrm{}\epsilon _r=\sigma \epsilon _1\mathrm{}\epsilon _r,\sigma S(n).$$
(3.7)
Hence the dimension of the subspace $`[S(n)]\epsilon _1\mathrm{}\epsilon _r`$ does not exceed the number of left cosets of $`S(n)`$ over the subgroup $`S(r)`$. Therefore, $`dim𝒜(n)`$ does not exceed
$$\underset{r=0}{\overset{n}{}}\left(\genfrac{}{}{0pt}{}{n}{r}\right)\frac{n!}{r!}$$
(3.8)
which coincides with (3.5).
Using Proposition 3.1 we shall sometimes identify $`A(n)`$ with the algebra $`𝒜(n)`$.
###### Corollary 3.2
The mapping
$$(k,k+1)(k,k+1),\epsilon _k0$$
(3.9)
defines an algebra homomorphism $`A(n)[S(n)]`$ which is identical on the subalgebra $`[S(n)]`$.
We shall call (3.9) the retraction homomorphism. It can be equivalently defined as follows. For any $`\gamma \mathrm{\Gamma }(n)`$, define its rank, denoted by $`\mathrm{rank}\gamma `$, as the number of $`1`$’s in the $`(0,1)`$-matrix representing $`\gamma `$. That is,
$$\mathrm{rank}\gamma =|\mathrm{dom}\gamma |=|\mathrm{range}\gamma |.$$
(3.10)
The rank of an element $`a=a_\gamma \gamma A(n)`$ is defined as the maximum of the ranks $`\mathrm{rank}\gamma `$ with $`a_\gamma 0`$. Now (3.9) can also be defined by setting for $`\gamma \mathrm{\Gamma }(n)`$
$$\gamma \{\begin{array}{cc}\gamma \hfill & \text{if}\mathrm{rank}\gamma =n,\hfill \\ 0\hfill & \text{if}\mathrm{rank}\gamma <n,\hfill \end{array}$$
(3.11)
and extending this to $`A(n)`$ by linearity.
For any $`0mn`$ denote by $`\mathrm{\Gamma }_m(n)`$ the subsemigroup of $`\mathrm{\Gamma }(n)`$ which consists of the matrices with first $`m`$ diagonal entries equal to $`1`$. Set $`A_m(n)=A(n)^{\mathrm{\Gamma }_m(n)}`$, the centralizer of $`\mathrm{\Gamma }_m(n)`$ in the algebra $`A(n)`$. In particular, $`A_0(n)`$ is the center of $`A(n)`$.
We extend $`\theta _n`$ to a linear map $`A(n)A(n1)`$.
###### Proposition 3.3
The restriction of $`\theta _n`$ to $`A_{n1}(n)A(n)`$ defines a unital algebra homomorphism
$$\theta _n:A_{n1}(n)A(n1).$$
(3.12)
Moreover,
$$\theta _n(A_m(n))A_m(n1)\text{for }m=0,1,\mathrm{},n1\text{.}$$
(3.13)
Proof. Note that any $`aA_{n1}(n)`$ commutes with $`\epsilon _n`$ because $`\epsilon _n`$ is contained in $`\mathrm{\Gamma }_{n1}(n)`$, and $`A_{n1}(n)`$ is the centralizer of $`\mathrm{\Gamma }_{n1}(n)`$ in $`A(n)`$. Since $`\epsilon _n`$ is an idempotent we have
$$\theta _n(a)=a\epsilon _n=\epsilon _na=\epsilon _na\epsilon _n,aA_{n1}(n).$$
(3.14)
Now, if $`a,bA_{n1}(n)`$ then
$$\theta _n(ab)=\epsilon _nab\epsilon _n=\epsilon _na\epsilon _n\epsilon _nb\epsilon _n=\theta _n(a)\theta _n(b).$$
(3.15)
It is clear that (3.12) preserves the unity.
Finally, we need to show that if $`aA_m(n)`$ and $`b\mathrm{\Gamma }_m(n1)`$ then $`\theta _n(a)`$ and $`b`$ commute. Regard $`b`$ as an element of $`\mathrm{\Gamma }(n)`$; then it lies in $`\mathrm{\Gamma }_m(n)`$ and commutes with $`\epsilon _n`$. This implies that $`a\epsilon _n`$ and $`b`$ commute, and so do $`\theta _n(a)`$ and $`b`$.
For $`\gamma \mathrm{\Gamma }(n)`$, set
$$\begin{array}{cc}\hfill J(\gamma )& =\{i|1in,\gamma _{ii}=0\},\hfill \\ \hfill \mathrm{deg}\gamma & =|J(\gamma )|.\hfill \end{array}$$
(3.16)
###### Proposition 3.4
Let $`\gamma ,\gamma ^{}\mathrm{\Gamma }(n)`$. Then
$$\mathrm{deg}\gamma \gamma ^{}\mathrm{deg}\gamma +\mathrm{deg}\gamma ^{}.$$
(3.17)
Moreover, the equality in (3.17) implies $`\gamma \gamma ^{}=\gamma ^{}\gamma `$.
Proof. Regard $`\gamma `$ and $`\gamma ^{}`$ as partial bijections of the set $`_n=\{1,\mathrm{},n\}`$. If $`\delta \mathrm{\Gamma }(n)`$ then $`_nJ(\delta )`$ is the set of $`\delta `$-invariant elements in the domain of $`\delta `$. This implies
$$\left(_nJ(\gamma )\right)\left(_nJ(\gamma ^{})\right)\left(_nJ(\gamma \gamma ^{})\right).$$
(3.18)
Therefore,
$$J(\gamma )J(\gamma ^{})J(\gamma \gamma ^{}),$$
(3.19)
and (3.17) follows. Finally, the equality in (3.17) implies that $`J(\gamma )`$ and $`J(\gamma ^{})`$ are disjoint. But then $`\gamma `$ and $`\gamma ^{}`$ must commute.
Using (3.16), we define a filtration of the space $`A(n)`$:
$$=A^0(n)A^1(n)\mathrm{}A^n(n)=A(n)$$
(3.20)
where $`A^M(n)`$ is spanned by the subset $`\{\gamma |\mathrm{deg}\gamma M\}\mathrm{\Gamma }(n)`$. By Proposition 3.4 this filtration is compatible with the algebra structure of $`A(n)`$, and the corresponding graded algebra
$$\mathrm{gr}A(n)=\underset{M}{}\left(A^M(n)/A^{M1}(n)\right)$$
(3.21)
is commutative. Note that for $`\gamma \mathrm{\Gamma }(n)`$ the degree $`\mathrm{deg}\theta _n(\gamma )`$ can be equal either to $`\mathrm{deg}\gamma `$ or $`\mathrm{deg}\gamma 1`$. Therefore the homomorphisms (3.13) are compatible with the filtration on $`A(n)`$.
###### Definition 3.5
For $`m=0,1,2,\mathrm{}`$ let $`A_m`$ be the projective limit of the infinite sequence
$$\mathrm{}A_m(n)\stackrel{\theta _n}{}A_m(n1)\mathrm{}A_m(m+1)\stackrel{\theta _{m+1}}{}A_m(m)$$
(3.22)
taken in the category of filtered algebras.
By the definition, an element $`aA_m`$ is a sequence $`(a_n|nm)`$ such that
$$a_nA_m(n),\theta _n(a_n)=a_{n1},\mathrm{deg}a:=\underset{nm}{sup}\mathrm{deg}a_n<\mathrm{}$$
(3.23)
with the componentwise operations. For $`nm`$ we shall denote by $`\theta ^{(n)}`$ the projection $`A_mA_m(n)`$ such that
$$\theta ^{(n)}(a)=a_n.$$
(3.24)
The $`M`$-th term of the filtered algebra $`A_m`$ will be denoted by $`A_m^M`$.
There are natural algebra homomorphisms $`A_mA_{m+1}`$ defined by
$$(a_n|nm)(a_n|nm+1)$$
(3.25)
where we use the inclusions $`A_m(n)A_{m+1}(n)`$ for $`n>m`$. These homomorphisms are injective because $`a_m`$ is uniquely determined by $`a_{m+1}`$.
###### Definition 3.6
The algebra $`A`$ is defined as the inductive limit (the union) of the algebras $`A_m`$ taken with respect to the embeddings $`A_mA_{m+1}`$, $`m0`$, defined in (3.25).
Since these embeddings preserve the filtration, $`A`$ is a filtered algebra. We will denote by $`A^M`$ the $`M`$-th term of the filtration, so that
$$A^M=\underset{m0}{}A_m^M.$$
(3.26)
###### Proposition 3.7
There exists a natural embedding $`A(\mathrm{})A`$ whose image consists of stable sequences $`a=(a_n)A`$.
Proof. Let $`bA(\mathrm{})`$. There exists $`m`$ such that $`bA(m)`$. Note that $`bA_m(n)`$ for any $`nm`$ since $`A(m)`$ and $`\mathrm{\Gamma }_m(n)`$ commute. Set
$$a=(a_n|nm)A_mA\text{with }a_nb\text{.}$$
(3.27)
The sequence $`aA`$ only depends on $`b`$ and not on the choice of $`m`$. The mapping $`ba`$ is clearly an algebra embedding.
###### Corollary 3.8
There is a natural algebra embedding $`[S(\mathrm{})]A`$.
###### Proposition 3.9
The center of the algebra $`A`$ coincides with $`A_0`$.
Proof. Recall that $`A_0(n)`$ is the center of $`A(n)`$. The subalgebra $`A_0`$ is contained in the center of $`A`$ since the sequences $`a=(a_n)A`$ are multiplied componentwise.
Conversely, if $`a`$ belongs to the center of $`A`$ then $`a`$ commutes with the subalgebra $`A(\mathrm{})A`$. This implies that for any $`n`$ the element $`a_n`$ is contained in $`A_0(n)`$, and so $`aA_0`$.
Remark. The same argument shows that the subalgebra $`A_mA`$ coincides with the centralizer in $`A`$ of the subalgebra
$$\underset{nm}{}[\mathrm{\Gamma }_m(n)]A(\mathrm{})A.$$
(3.28)
Note that the centers of both algebras $`[S(\mathrm{})]`$ and $`A(\mathrm{})`$ are trivial. However, as it will be shown in the next section, the center $`A_0`$ of the algebra $`A`$ has a rich structure.
###### Proposition 3.10
For any tame representation $`T`$ of the group $`S(\mathrm{})`$, the subspace $`H_{\mathrm{}}(T)H(T)`$ admits a natural structure of an $`A`$-module such that for any $`m`$ the subspace $`H_m(T)`$ is invariant with respect to the subalgebra $`A_m`$ (and hence $`H_n(T)`$ is invariant with respect to $`A_m`$ for $`nm`$).
Proof. Let $`aA`$ and $`hH_{\mathrm{}}(T)`$. Choose $`m`$ such that $`aA_m`$. Then we may write $`a=(a_n|nm)`$. Let us prove that
$$a_nh=a_{n+1}h,nm.$$
(3.29)
Consider the family of representations $`\{𝒯_n\}`$ associated with $`T`$, which has been introduced in Section 2.2. Each $`𝒯_n`$ is a representation of the semigroup $`\mathrm{\Gamma }(n)`$ in the space $`H_n(T)`$ and so it can be extended to a representation of the semigroup algebra $`A(n)`$ in the same space. Recall that $`𝒯_{n+1}(\epsilon _{n+1})`$ projects $`H_{n+1}(T)`$ onto its subspace $`H_n(T)`$. Since $`h`$ is already contained in $`H_n(T)`$ (as we assume $`nm`$), we have $`𝒯_{n+1}(\epsilon _{n+1})h=h`$, so that
$$T_{n+1}(1\epsilon _{n+1})h=0.$$
(3.30)
This implies that $`h`$ is annihilated by the left ideal $`I(n+1)A(n+1)`$. Since $`a_na_{n+1}I(n+1)`$, this implies (3.29).
Define a mapping
$$A\times H_{\mathrm{}}(T)H_{\mathrm{}}(T),(a,h)a_mh$$
(3.31)
where $`m`$ is so large that $`aA_m`$ and $`hH_m(T)`$. Note that under this assumption $`ahH_m(T)`$.
The mapping (3.31) is clearly bilinear and $`1h=h`$. The multiplicativity property $`(ab)h=a(bh)`$ follows from the definition of the multiplication in $`A_m`$.
###### Proposition 3.11
If $`T`$ is an irreducible tame representation of $`S(\mathrm{})`$ then $`H_{\mathrm{}}(T)`$ is irreducible as an $`A`$-module. In particular, the center $`A_0`$ of $`A`$ acts by scalar operators.
Proof. The first claim is obvious because $`H_{\mathrm{}}(T)`$ is already irreducible as a $`A(\mathrm{})`$-module (see Proposition 3.7). To prove the second claim consider an element $`aA_0`$ as an operator in $`H_{\mathrm{}}(T)`$. It suffices to show that $`a`$ has an eigenvalue. The result will then follow by a standard argument using Schur’s lemma.
Assume that $`a`$ has no eigenvalues. Let us show first that $`a`$ is algebraically independent over $``$. Indeed, let $`P(x)[x]`$ be a nonzero polynomial of a minimum degree such that $`P(a)=0`$. Then $`P(x)=(x\alpha )Q(x)`$ for a certain $`\alpha `$ and a polynomial $`Q(x)[x]`$. Since $`Q(a)`$ is a nonzero operator, there is a vector $`vH_{\mathrm{}}(T)`$ such that $`w:=Q(a)v0`$. Then $`w`$ is an eigenvector for $`a`$ with the eigenvalue $`\alpha `$. Contradiction.
We note now that the space $`H_{\mathrm{}}(T)`$ has countable dimension and then use a version of Dixmier’s argument as follows.
Since $`a`$ is algebraically independent over $``$, the field $`(a)`$ is embedded in the endomorphism algebra of $`H_{\mathrm{}}(T)`$. This implies that the dimension of $`H_{\mathrm{}}(T)`$ is at least as large as the dimension of $`(a)`$ over $``$, but the latter is continuum. This contradiction completes the proof.
## 4 The structure of the algebra $`A_0`$
In the last two sections we aim to describe the structure of the algebras $`A_m`$. Here we consider the commutative algebra $`A_0`$; see Proposition 3.9. We construct generators of $`A_0`$ and show that it is isomorphic to the algebra of shifted symmetric functions.
### 4.1 Generators of $`A_0`$
Let $`Z(S(n))`$ denote the center of the algebra $`[S(n)]`$. For $`0Mn`$ denote by $`Z^M(S(n))`$ the $`M`$-th term of the filtration on $`Z(S(n))`$ inherited from the algebra $`A(n)`$; see (3.20). Note that
$$Z^0(S(n))=Z^1(S(n))=\mathrm{\hspace{0.17em}1}$$
(4.1)
because, for a permutation $`sS(n)`$ the inequality $`\mathrm{deg}s1`$ implies $`s=1`$.
For any partition $`=(M_1,\mathrm{},M_r)`$ with
$$||=M_1+\mathrm{}+M_rn$$
(4.2)
introduce the element $`c_n^{}`$ of the group algebra $`[S(n)]`$ as follows
$$c_n^{}=(i_1,\mathrm{},i_{M_1})(j_1,\mathrm{},j_{M_2})\mathrm{}(k_1,\mathrm{},k_{M_r})$$
(4.3)
where the sum is taken over the sequences $`i_1,\mathrm{},i_{M_1};j_1,\mathrm{},j_{M_2};\mathrm{};k_1,\mathrm{},k_{M_r}`$ of $`||`$ pairwise distinct indices taken from $`_n`$. By $`(i_1,\mathrm{},i_{M_1})`$ etc. in (4.3) we denote cycles in the symmetric group. For the empty partition $`\mathrm{}`$ we set $`c_n^{\mathrm{}}=1`$. Note that $`c_n^{(1)}=_{i=1}^n(i)=n\mathrm{\hspace{0.17em}1}.`$ Given two partitions $``$ and $``$ we denote by $``$ the partition whose parts are those of $``$ and $``$ rewritten in the decreasing order. We have
$$c_n^{1\mathrm{}1}=(n||)\mathrm{}(n||p+1)c_n^{}$$
(4.4)
where $`p`$ stands for the number of $`1`$’s in the left hand side of the relation.
By definition (3.16) of the degree of an element of $`\mathrm{\Gamma }(n)`$ we have
$$\mathrm{deg}c_n^{(1)}=0,\text{and}\mathrm{deg}c_n^{(M)}=M\text{for }M2\text{.}$$
(4.5)
More generally,
$$\mathrm{deg}c_n^{(M_1,\mathrm{},M_r)}=\underset{i,M_i2}{}M_i.$$
(4.6)
###### Proposition 4.1
Each of the families
$$c_n^{},||=n,$$
(4.7)
and
$$c_n^{},||n\text{and }\text{ has no part equal to }1\text{,}$$
(4.8)
forms a basis of $`Z(S(n))`$. Moreover, the elements of degree $`M`$ of each family form a basis of $`Z^M(S(n))`$.
Proof. The elements (4.7) are proportional to the characteristic functions of the conjugacy classes of the group $`S(n)`$ and so, they form a basis of $`Z(S(n))`$. By (4.4) the elements of type (4.8) are proportional to those of type (4.7).
###### Proposition 4.2
For any two partitions $`=(M_1,\mathrm{},M_r)`$ and $`=(L_1,\mathrm{},L_t)`$ with $`||+||n`$ we have
$$c_n^{}c_n^{}=c_n^{}+(\mathrm{}),$$
(4.9)
where $`(\mathrm{})`$ stands for a linear combination of the elements $`c_n^𝒦`$ with $`|𝒦|<||+||`$.
Proof. For a permutation $`sS(n)`$ or $`sS(\mathrm{})`$ define its support as
$$\mathrm{supp}s=\{i_n|s(i)i\}\text{or}\mathrm{supp}s=\{i|s(i)i\},$$
(4.10)
respectively. (The degree of a permutation is then given by $`\mathrm{deg}s=|\mathrm{supp}s|`$; cf. (3.16)).
Let $`sS(n)`$ be a permutation which occurs in the expansion of $`c_n^{}`$, that is, $`s`$ is of cycle type $`1\mathrm{}1`$ (with $`n||`$ units). Similarly, let $`s^{}`$ be a permutation occurring in $`c_n^{}`$. If the supports $`\mathrm{supp}s`$ and $`\mathrm{supp}s^{}`$ are disjoint then $`s`$ and $`s^{}`$ commute, and the product $`ss^{}`$ occurs in the expansion of $`c_n^{}`$. In particular, $`\mathrm{deg}ss^{}=||+||`$. If $`\mathrm{supp}s`$ and $`\mathrm{supp}s^{}`$ have a non-empty intersection then the degree of $`ss^{}`$ is strictly less than $`||+||`$.
Remark. A detailed investigation of the structure constants for the products of type (4.9) have been recently given by Ivanov and Kerov .
###### Corollary 4.3
Let $`k=(k_1,\mathrm{},k_n)`$ run over the $`n`$-tuples of non-negative integers such that $`2k_2+\mathrm{}+nk_nn`$. Then the monomials
$$(c_n^{(2)})^{k_2}\mathrm{}(c_n^{(n)})^{k_n}$$
(4.11)
form a basis of $`Z(S(n))`$. Moreover, for any $`M0`$, the elements (4.11) with $`2k_2+\mathrm{}+nk_nM`$ form a basis of $`Z^M(S(n))`$.
Proof. It suffices to prove that
$$(c_n^{(2)})^{k_2}\mathrm{}(c_n^{(n)})^{k_n}=c_n^{}+(\mathrm{})$$
(4.12)
where $`=2^{k_2}\mathrm{}n^{k_n}`$ and $`(\mathrm{})`$ stands for a certain linear combination of the elements $`c_n^{^{}}`$ with $`|^{}|<||=2k_2+\mathrm{}+nk_n`$. But this follows from Proposition 4.2.
Now we will define analogs of the elements $`c_n^{}`$ for the algebra $`A_0(n)`$. Namely, for any partition $`=(M_1,\mathrm{},M_r)`$ with $`||n`$ set
$$\mathrm{\Delta }_n^{}=(i_1,\mathrm{},i_{M_1})(j_1,\mathrm{},j_{M_2})\mathrm{}(k_1,\mathrm{},k_{M_r})(1\epsilon _{i_1})\mathrm{}(1\epsilon _{k_{M_r}})$$
(4.13)
where, as in (4.3), the sum is taken over all sequences of $`||`$ pairwise distinct indices taken from $`_n`$. In particular,
$$\mathrm{\Delta }_n^{(1)}=\underset{i=1}{\overset{n}{}}(1\epsilon _i).$$
(4.14)
For the empty partition $`\mathrm{}`$ we set $`\mathrm{\Delta }_n^{\mathrm{}}=1`$. By (3.16), we have
$$\mathrm{deg}\mathrm{\Delta }_n^{}=||\text{for any partition },$$
(4.15)
cf. (4.6). Note that $`\mathrm{\Delta }_n^{}`$ can also be written as
$$\mathrm{\Delta }_n^{}=(1\epsilon _{i_1})\mathrm{}(1\epsilon _{k_{M_r}})(i_1,\mathrm{},i_{M_1})(j_1,\mathrm{},j_{M_2})\mathrm{}(k_1,\mathrm{},k_{M_r}),$$
(4.16)
and as
$$\begin{array}{cc}\hfill \mathrm{\Delta }_n^{}=(1\epsilon _{i_1})\mathrm{}(1\epsilon _{k_{M_r}})& (i_1,\mathrm{},i_{M_1})(j_1,\mathrm{},j_{M_2})\hfill \\ \hfill \mathrm{}& (k_1,\mathrm{},k_{M_r})(1\epsilon _{i_1})\mathrm{}(1\epsilon _{k_{M_r}}).\hfill \end{array}$$
(4.17)
Indeed, $`(1\epsilon _{i_1})\mathrm{}(1\epsilon _{k_{M_r}})`$ is invariant under the conjugation by the permutation $`(i_1,\mathrm{},i_{M_1})(j_1,\mathrm{},j_{M_2})\mathrm{}(k_1,\mathrm{},k_{M_r})`$ which implies (4.16). To derive (4.17) it suffices to note that $`(1\epsilon _{i_1})\mathrm{}(1\epsilon _{k_r})`$ is an idempotent.
###### Proposition 4.4
The element $`\mathrm{\Delta }_n^{}`$ belongs to $`A_0(n)`$ for any $``$.
Proof. Since $`\mathrm{\Gamma }(n)`$ is generated by the group $`S(n)`$ and the pairwise commuting idempotents $`\epsilon _1,\mathrm{},\epsilon _n`$, it suffices to show that $`\mathrm{\Delta }_n^{}`$ commutes both with $`S(n)`$ and with the $`\epsilon _i`$. The first claim is clear since $`\mathrm{\Delta }_n^{}`$ is invariant under the conjugation by the elements of $`S(n)`$. To prove the second claim, we observe that any $`\epsilon _l`$, $`1ln`$, commutes with any term
$$\sigma =(i_1,\mathrm{},i_{M_1})(j_1,\mathrm{},j_{M_2})\mathrm{}(k_1,\mathrm{},k_{M_r})(1\epsilon _{i_1})\mathrm{}(1\epsilon _{k_{M_r}})$$
(4.18)
in (4.13). Indeed, this is clear if $`l`$ does not occur in the set of indices in (4.18) because $`\epsilon _l`$ commutes with the corresponding cycle. But if $`l`$ coincides with one of the indices $`i_1,\mathrm{},k_{M_r}`$, then $`\epsilon _l\sigma =\sigma \epsilon _l=0`$. This follows from (4.17) and the relation $`(1\epsilon _l)\epsilon _l=0`$.
###### Proposition 4.5
We have
$$\theta _n(\mathrm{\Delta }_n^{})=\mathrm{\Delta }_{n1}^{}$$
(4.19)
where we adopt the convention that
$$\mathrm{\Delta }_k^{}=0\text{if}||>k.$$
(4.20)
Proof. By the definition of the projection $`\theta _n`$ (see Section 3) we need to calculate $`\mathrm{\Delta }_n^{}\epsilon _n`$. However, as it follows from the proof of Proposition 4.4, the effect of multiplying $`\mathrm{\Delta }_n^{}`$ by $`\epsilon _n`$ reduces to striking from (4.13) all terms (4.18) such that $`n`$ occurs among the corresponding indices. If $`||=n`$, then all the terms are vanished, so that the result of the multiplication is $`0`$. If $`||<n`$, then the terms that survive are just the terms of the sum defining $`\mathrm{\Delta }_{n1}^{}`$.
###### Corollary 4.6
For any partition $``$, there exists an element $`\mathrm{\Delta }^{}A_0`$ such that
$$\theta ^{(n)}(\mathrm{\Delta }^{})=\mathrm{\Delta }_n^{}\text{for any }n1$$
(4.21)
with the convention (4.20).
Proof. By Proposition 4.4, $`\mathrm{\Delta }_n^{}A_0(n)`$. Now we apply Proposition 4.5 and note that the degrees of the elements $`\mathrm{\Delta }_n^{}`$ are uniformly bounded by (4.15).
We now aim to prove an analog of Proposition 4.1 for the algebra $`A_0(n)`$; see Proposition 4.10 below. For this we need the following three lemmas.
Let $`I(n)=A(n)(1\epsilon _n)`$ denote the left ideal of the algebra $`A(n)`$ generated by the element $`1\epsilon _n`$.
###### Lemma 4.7
For any $`n`$,
$$I(n)A_0(n)=Z(S(n))(1\epsilon _1)\mathrm{}(1\epsilon _n).$$
(4.22)
Proof. First suppose that $`xA(n)`$ can be written as $`y(1\epsilon _1)\mathrm{}(1\epsilon _n)`$ where $`yZ(S(n))`$. The argument of the proof of Proposition 4.4 shows that $`xA_0(n)`$. Moreover, we obviously have $`xI(n)`$.
Conversely, suppose $`xI(n)A_0(n)`$. Then $`x\epsilon _n=0`$. Using the invariance of $`x`$ under the conjugation by elements of $`S(n)`$ we also obtain $`x\epsilon _i=0`$ for $`i=1,\mathrm{},n`$. Therefore $`x`$ is invariant under the right multiplication by $`(1\epsilon _1)\mathrm{}(1\epsilon _n)`$. Further, we may write
$$x=y+\underset{r=1}{\overset{n}{}}\underset{1i_1<\mathrm{}<i_rn}{}y_{i_1\mathrm{}i_r}\epsilon _{i_1}\mathrm{}\epsilon _{i_r},$$
(4.23)
where $`y`$ and all the $`y_{i_1\mathrm{}i_r}`$ are elements of $`[S(n)]`$; see Proposition 3.1. Multiplying this relation by $`(1\epsilon _1)\mathrm{}(1\epsilon _n)`$ on the right we obtain
$$x=y(1\epsilon _1)\mathrm{}(1\epsilon _n).$$
(4.24)
Finally, for any $`sS(n)`$ we may write
$$x=sxs^1=sys^1(1\epsilon _1)\mathrm{}(1\epsilon _n).$$
(4.25)
Averaging over $`sS(n)`$ turns $`y`$ into an element of $`Z(S(n))`$.
For a subset $`I=\{i_1,\mathrm{},i_k\}`$ in $`_n`$ we put $`\epsilon _I=\epsilon _{i_1}\mathrm{}\epsilon _{i_k}`$, and for $`sS(n)`$ set
$$Q(s)=\{i_n|s_{ii}=1\}=_n\mathrm{supp}s.$$
(4.26)
###### Lemma 4.8
The mapping
$$s\gamma ,\gamma =s\epsilon _{Q(s)}=\epsilon _{Q(s)}s,$$
(4.27)
defines a bijection of $`S(n)`$ onto the set of all $`\gamma \mathrm{\Gamma }(n)`$ satisfying the conditions
$`\mathrm{dom}\gamma `$ $`=\mathrm{range}\gamma ,`$ (4.28)
$`\mathrm{deg}\gamma `$ $`=n.`$ (4.29)
Proof. The effect of the multiplication of $`s`$ by $`\epsilon _{Q(s)}`$ from the left or from the right consists of replacing all the 1’s on the diagonal by zeros. This implies (4.28), and (4.29) is obvious.
Conversely, let $`\gamma \mathrm{\Gamma }(n)`$ satisfy (4.28) and (4.29). Relation (4.28) means that for any $`i_n`$ the $`i`$-th row and the $`i`$-th column are zero or non-zero at the same time, whereas (4.29) means that all the diagonal entries of $`\gamma `$ are zero. Now, let the matrix $`\sigma `$ be defined as follows. Set $`\sigma _{ii}=1`$ if the $`i`$-th row (and the $`i`$-th column) of $`\gamma `$ is zero, and set $`\sigma _{ij}=\gamma _{ij}`$ for $`ij`$. It is easy to see that $`\sigma S(n)`$ and that $`\gamma `$ is the image of $`\sigma `$ under the mapping (4.27).
###### Lemma 4.9
The restriction of the projection $`\theta _n:A_0(n)A_0(n1)`$ to the subspace $`A_0^{n1}(n)`$ is injective.
Proof. Let $`xA_0(n)`$ and $`\theta _n(x)=0`$. We will show that $`\mathrm{deg}x=n`$ unless $`x=0`$. By Lemma 4.7, $`x`$ can be written as a linear combination of the elements
$$s(1\epsilon _1)\mathrm{}(1\epsilon _n)=\underset{I_n}{}(1)^{|I|}s\epsilon _I,sS(n).$$
(4.30)
Rewrite this as
$$s(1\epsilon _1)\mathrm{}(1\epsilon _n)=\underset{IQ(s)}{}(1)^{|I|}s\epsilon _I+\underset{IQ(s)}{}(1)^{|I|}s\epsilon _I.$$
(4.31)
Then the terms of the first sum are of degree $`n`$ whereas those of the second sum are of degree strictly less than $`n`$. So it suffices to prove that the elements
$$\underset{IQ(s)}{}(1)^{|I|}s\epsilon _I,sS(n),$$
(4.32)
are linearly independent. Note that
$`\mathrm{rank}s\epsilon _I`$ $`=n|Q(s)|\text{if }I=Q(s)\text{,}`$ (4.33)
$`\mathrm{rank}s\epsilon _I`$ $`<n|Q(s)|\text{if }IQ(s);`$ (4.34)
see (3.10). Write $`S(n)`$ as the disjoint union of $`n+1`$ subsets:
$$S(n)=\underset{k=0}{\overset{n}{}}\{sS(n)|n|Q(s)|=k\}.$$
(4.35)
If $`s`$ belongs to the $`k`$-th subset then
$$\mathrm{rank}\left(\underset{IQ(s)}{}(1)^{|I|}s\epsilon _I\right)=k.$$
(4.36)
Moreover, only one term of the sum in (4.36) has rank $`k`$, namely that with $`I=Q(s)`$. Finally, it remains to note that by Lemma 4.8 all the elements $`s\epsilon _{Q(s)}`$ with $`sS(n)`$ are pairwise distinct elements of $`\mathrm{\Gamma }(n)`$.
###### Proposition 4.10
For any $`n`$ the elements $`\mathrm{\Delta }_n^{}`$, where $``$ is any partition with $`||n`$, form a basis of $`A_0(n)`$. Furthermore, for any $`M`$ such that $`0Mn`$ these elements with $`||M`$ form a basis of $`A_0^M(n)`$.
Proof. The first claim of the proposition will follow from the second one. We will prove the second claim using induction on $`n`$. The claim is obviously true for $`n=1`$. Assume that $`n2`$ and $`Mn1`$. By the induction hypothesis the elements $`\mathrm{\Delta }_{n1}^{}`$ with $`||M`$ form a basis of $`A_0^M(n1)`$. By Proposition 4.5 the image of $`\mathrm{\Delta }_n^{}`$ under $`\theta _n`$ is $`\mathrm{\Delta }_{n1}^{}`$. By Lemma 4.9 the restriction $`\theta _nA_0^M(n)`$ is injective. Therefore, the elements $`\mathrm{\Delta }_n^{}`$ with $`||M`$ form a basis in $`A_0^M(n)`$.
Further, let us show that the elements $`\mathrm{\Delta }_n^{}`$ with $`||=n`$ form a basis of $`I(n)A_0(n)`$. Note that
$$\mathrm{\Delta }_n^{}=c_n^{}(1\epsilon _1)\mathrm{}(1\epsilon _n).$$
(4.37)
By Proposition 4.1 the elements $`c_n^{}`$, where $``$ runs over the set of partitions of $`n`$, form a basis of $`Z(S(n))`$. Due to Lemma 4.7 it now remains to check that the elements $`c_n^{}`$, being multiplied by $`(1\epsilon _1)\mathrm{}(1\epsilon _n)`$, remain linearly independent. However, this follows from the fact that the composite map
$$[S(n)]A(n)[S(n)]$$
(4.38)
is the identity map; here the first arrow is the multiplication by $`(1\epsilon _1)\mathrm{}(1\epsilon _n)`$, and the second arrow is the retraction homomorphism (3.9).
Finally, let us show that
$$A_0(n)=A_0^{n1}(n)\left(I(n)A_0(n)\right).$$
(4.39)
Indeed, as it was shown above, $`\theta _n`$ maps $`A_0^{n1}(n)`$ onto $`A_0^{n1}(n1)=A_0(n1)`$ . Since $`I(n)A_0(n)`$ is the kernel of the restriction $`\theta _nA_0(n)`$ and since $`\theta _n(A_0(n))`$ is contained in $`A_0(n1)`$, we obtain the decomposition
$$A_0(n)=A_0^{n1}(n)+\left(I(n)A(n)\right).$$
(4.40)
Lemma 4.9 implies that
$$A_0^{n1}(n)I(n)=\{0\}$$
(4.41)
and (4.39) follows.
To complete the proof we need to show that the elements $`\mathrm{\Delta }_n^{}`$ with $`0||n`$ form a basis of $`A_0(n)`$. However, the elements with $`||<n`$ form a basis of the first component of the decomposition (4.39), whereas the elements with $`||=n`$ form a basis in the second component of this decomposition.
The following is an analog of Corollary 4.3.
###### Corollary 4.11
Let $`k=(k_1,\mathrm{},k_n)`$ run over the $`n`$-tuples of non-negative integers such that $`k_1+2k_2+\mathrm{}+nk_nn.`$ Then the monomials
$$(\mathrm{\Delta }_1^{(1)})^{k_1}\mathrm{}(\mathrm{\Delta }_n^{(n)})^{k_n}$$
(4.42)
form a basis of $`A_0(n)`$. Moreover, for any $`M0`$, the monomials (4.42) with $`k_1+2k_2+\mathrm{}+nk_nM`$ form a basis of $`A_0^M(n)`$.
Proof. It suffices to prove that
$$(\mathrm{\Delta }_1^{(1)})^{k_1}\mathrm{}(\mathrm{\Delta }_n^{(n)})^{k_n}=\mathrm{\Delta }_n^{}+(\mathrm{})$$
(4.43)
where $`=1^{k_1}2^{k_2}\mathrm{}n^{k_n}`$ and $`(\mathrm{})`$ stands for a linear combination of the elements $`\mathrm{\Delta }_n^{^{}}`$ with $`|^{}|<||`$. Then our claim will follow from Proposition 4.10. To prove (4.43) we verify that for any partitions $`=(M_1,\mathrm{},M_r)`$ and $`=(L_1,\mathrm{},L_t)`$ with $`||+||n`$
$$\mathrm{\Delta }_n^{}\mathrm{\Delta }_n^{}=\mathrm{\Delta }_n^{}+(\mathrm{}),$$
(4.44)
where the rest term $`(\mathrm{})`$ has degree strictly less than $`||+||`$ and so, it is a linear combination of elements $`\mathrm{\Delta }_n^𝒦`$ with $`|𝒦|<||+||`$. Write
$$\mathrm{\Delta }_n^{}=\delta _I,\mathrm{\Delta }_n^{}=\delta _J^{}.$$
(4.45)
Here $`I`$ is a sequence $`i_1,\mathrm{},i_{||}`$ of pairwise distinct indices taken from $`_n`$ and
$$\delta _I=(i_1,\mathrm{},i_{M_1})\mathrm{}(i_{M_1+\mathrm{}+M_{r1}+1},\mathrm{},i_{||})\underset{p=1}{\overset{||}{}}(1\epsilon _{i_p});$$
(4.46)
the $`\delta _J^{}`$ are the corresponding elements for the partition $``$. Then
$$\mathrm{\Delta }_n^{}\mathrm{\Delta }_n^{}=\underset{I,J}{}\delta _I\delta _J^{}=\underset{IJ=\mathrm{}}{}\delta _I\delta _J^{}+\underset{IJ\mathrm{}}{}\delta _I\delta _J^{}.$$
(4.47)
The first sum on the right hand side of (4.47) is $`\mathrm{\Delta }_n^{}`$ whereas the second sum is of degree strictly less than $`||+||`$.
Consider the elements $`\mathrm{\Delta }^{}A_0`$ introduced in Corollary 4.6. We shall denote by $``$ the set of all partitions.
###### Theorem 4.12
The elements $`\mathrm{\Delta }^{}`$, $``$ form a basis of the algebra $`A_0`$. Moreover, for any $`M0`$, the elements $`\mathrm{\Delta }^{}`$ with $`||M`$ form a basis of the $`M`$-th subspace $`A_0^M`$ in $`A_0`$.
Proof. The first claim follows from the second one. The second claim follows from Proposition 4.10 and the definition of $`A_0^M`$ as the projective limit of the spaces $`A_0^M(n)`$.
###### Corollary 4.13
For $`n>M`$, the mapping
$$\theta _n:A_0^M(n)A_0^M(n1)$$
(4.48)
is an isomorphism of vector spaces and so is the mapping
$$\theta ^{(n)}:A_0^MA_0^M(n),nM.$$
(4.49)
In particular, $`dimA_0^M<\mathrm{}`$.
###### Theorem 4.14
The monomials
$$(\mathrm{\Delta }^{(1)})^{k_1}(\mathrm{\Delta }^{(2)})^{k_2}\mathrm{}$$
(4.50)
with $`k_1,k_2,\mathrm{}_+`$ and $`k_1+2k_2+\mathrm{}<\mathrm{}`$ form a basis of the algebra $`A_0`$. Moreover, for any $`M0`$ the monomials (4.50) with $`k_1+2k_2+\mathrm{}M`$ form a basis of the subspace $`A_0^M`$.
Proof. It suffices to check that
$$(\mathrm{\Delta }^{(1)})^{k_1}(\mathrm{\Delta }^{(2)})^{k_2}\mathrm{}\mathrm{\Delta }^{}modA_0^{M1}$$
(4.51)
where $`=1^{k_1}2^{k_2}\mathrm{}`$ and $`M=||`$. However, this follows from the relation (4.43).
###### Corollary 4.15
The elements $`\mathrm{\Delta }^{(1)},\mathrm{\Delta }^{(2)},\mathrm{}`$ are algebraically independent and generate the algebra $`A_0`$.
### 4.2 The algebra $`\mathrm{\Lambda }^{}`$ of shifted symmetric functions, and the isomorphism $`A_0\mathrm{\Lambda }^{}`$
Let $`\mathrm{\Lambda }^{}(n)[x_1,\mathrm{},x_n]`$ denote the subalgebra of polynomials in $`n`$ variables $`x_1,\mathrm{},x_n`$ which are symmetric in the new variables
$$y_1=x_11,y_2=x_22,\mathrm{},y_n=x_nn.$$
(4.52)
Following , we refer to $`\mathrm{\Lambda }^{}(n)`$ as the algebra of shifted symmetric polynomials in $`n`$ variables. We equip $`\mathrm{\Lambda }^{}(n)`$ with the filtration with respect to the usual degree of polynomials. Set $`\mathrm{\Lambda }^{}(0)=`$ and for $`n1`$ define the projection $`\mathrm{\Lambda }^{}(n)\mathrm{\Lambda }^{}(n1)`$ by specializing $`x_n=0`$. Note that this projection preserves the filtration.
###### Definition 4.16
The algebra $`\mathrm{\Lambda }^{}`$ of shifted symmetric functions is the projective limit of the filtered algebras $`\mathrm{\Lambda }^{}(n)`$ as $`n\mathrm{}`$.
In other words, an element $`f\mathrm{\Lambda }^{}`$ is a sequence $`(f_n|n0)`$ such that
(i) $`f_n\mathrm{\Lambda }^{}(n)`$ for any $`n`$;
(ii) for any $`n1`$, $`f_nf_{n1}`$ under the projection $`\mathrm{\Lambda }^{}(n)\mathrm{\Lambda }^{}(n1)`$;
(iii) $`\mathrm{deg}f_n`$ remains bounded as $`n\mathrm{}`$.
For an element $`f=(f_n)\mathrm{\Lambda }^{}`$, we define its degree by
$$\mathrm{deg}f=\underset{n}{sup}\mathrm{deg}f_n,$$
(4.53)
and for $`M=0,1,\mathrm{}`$ we denote by $`(\mathrm{\Lambda }^{})^M`$ the subspace in $`\mathrm{\Lambda }^{}`$ consisting of the elements of degree $`M`$. The algebra $`\mathrm{\Lambda }^{}`$ was first introduced in . A detailed study of $`\mathrm{\Lambda }^{}`$ is contained in .
Note an evident similarity between the shifted symmetric functions and the symmetric functions. Recall (see ) that the algebra $`\mathrm{\Lambda }`$ of symmetric functions is defined as the projective limit as $`n\mathrm{}`$ of the graded algebras $`\mathrm{\Lambda }(n)[x_1,\mathrm{},x_n]`$ of symmetric polynomials in $`n`$ variables. A difference between $`\mathrm{\Lambda }^{}`$ and $`\mathrm{\Lambda }`$ consists in a shift of variables and the replacement of the gradation by a filtration. The algebra $`\mathrm{\Lambda }^{}`$ may be viewed as a deformation of the algebra $`\mathrm{\Lambda }`$. Indeed, let $`h`$ be a numerical parameter, and let $`\mathrm{\Lambda }_h^{}`$ be defined similarly to $`\mathrm{\Lambda }^{}`$ but with $`y_i=x_iih`$ instead of (4.52). Then the algebras $`\mathrm{\Lambda }_h^{}`$ with $`h0`$, are naturally isomorphic to each other. Moreover, $`\mathrm{\Lambda }_1^{}`$ coincides with $`\mathrm{\Lambda }^{}`$ while $`\mathrm{\Lambda }_0^{}`$ coincides with $`\mathrm{\Lambda }`$. Another relation between $`\mathrm{\Lambda }^{}`$ and $`\mathrm{\Lambda }`$ is given by
###### Proposition 4.17
The graded algebra
$$\mathrm{gr}\mathrm{\Lambda }^{}=\underset{M=1}{\overset{\mathrm{}}{}}\left((\mathrm{\Lambda }^{})^M/(\mathrm{\Lambda }^{})^{M1}\right)$$
(4.54)
is isomorphic to the algebra $`\mathrm{\Lambda }`$.
Proof. For any $`M1`$ and any $`n`$, $`\mathrm{\Lambda }^{}(n)^M/\mathrm{\Lambda }^{}(n)^{M1}`$ is naturally isomorphic to the $`M`$-th homogeneous component of the algebra $`\mathrm{\Lambda }(n)[x_1,\mathrm{},x_n]`$. Moreover, this isomorphism is compatible with the projections $`\mathrm{\Lambda }^{}(n)\mathrm{\Lambda }^{}(n1)`$ and $`\mathrm{\Lambda }(n)\mathrm{\Lambda }(n1)`$. This yields an isomorphism $`\mathrm{gr}\mathrm{\Lambda }^{}\mathrm{\Lambda }`$.
In the following example we give some families of generators of the algebra $`\mathrm{\Lambda }^{}`$. Note that there also exist other important families analogous to the basic symmetric functions; see .
###### Example 4.18
For $`M=1,2,\mathrm{}`$, elements $`e_M`$, $`h_M`$, and $`p_M`$ defined by the formulas below, are shifted symmetric functions:
$`E(t)`$ $`=1+{\displaystyle \underset{M=1}{\overset{\mathrm{}}{}}}e_Mt^M={\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1+(x_kk)t}{1kt}},`$
$`H(t)`$ $`=1+{\displaystyle \underset{M=1}{\overset{\mathrm{}}{}}}h_Mt^M={\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1+kt}{1(x_kk)t}},`$
$`p_M`$ $`={\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}\left((x_kk)^M(k)^M\right).`$
The generating functions satisfy the following relations; cf. :
$$E(t)H(t)=1,\underset{k=1}{\overset{\mathrm{}}{}}p_Mt^M=t\frac{d}{dt}\mathrm{log}H(t).$$
(4.55)
###### Proposition 4.19
The algebra $`\mathrm{\Lambda }^{}`$ is isomorphic to the algebra of polynomials in countably many generators. Furthermore, we have
$$\mathrm{\Lambda }^{}=[e_1,e_2,\mathrm{}]=[h_1,h_2,\mathrm{}]=[p_1,p_2,\mathrm{}].$$
(4.56)
Proof. The corresponding statement for the algebra $`\mathrm{\Lambda }`$ of symmetric functions is well known, see \[13, Ch. 1, Section 2\]. Now, we apply Proposition 4.17 and note that the image of the shifted symmetric function $`e_M`$, $`h_M`$ or $`p_M`$ in the space $`(\mathrm{\Lambda }^{})^M/(\mathrm{\Lambda }^{})^{M1}\mathrm{\Lambda }^M`$ is the corresponding $`M`$-th symmetric function (elementary, complete or power sum). This implies that each of the three families is algebraically independent and generates the algebra $`\mathrm{\Lambda }^{}`$.
Let $`\mathrm{Fun}`$ denote the algebra of complex functions on the set of partitions $``$. By Propositions 3.10 and 3.11, there is an algebra homomorphism
$$A_0\mathrm{Fun},a\widehat{a},$$
(4.57)
such that for $`aA_0`$ and $`\lambda `$, the element $`a`$ acts in $`H_{\mathrm{}}(T_\lambda )`$ as the scalar operator $`\widehat{a}(\lambda )1`$. On the other hand, any $`\lambda `$ can be viewed as a sequence $`(\lambda _1,\lambda _2,\mathrm{},0,0,\mathrm{})`$ with finitely many non-zero coordinates, and so, any element of $`\mathrm{\Lambda }^{}`$ may be viewed as a function on $``$. Thus we obtain an algebra homomorphism $`\mathrm{\Lambda }^{}\mathrm{Fun}`$ which is clearly an embedding.
Let $`\lambda `$ be a partition with $`m=|\lambda |n`$. Consider the corresponding irreducible representation $`\pi _\lambda `$ of $`S(m)`$, and the representation $`𝒯_n(\lambda )`$ of the semigroup $`\mathrm{\Gamma }(n)`$; see Section 2.2.
###### Proposition 4.20
The eigenvalue of the central element $`\mathrm{\Delta }_n^{(r)}`$ in $`𝒯_n(\lambda )`$ is $`0`$ if $`r>m`$. If $`rm`$ then the eigenvalue coincides with that of the element $`c_m^{(r)}`$ in the representation $`\pi _\lambda `$ of $`S(m)`$.
Proof. Recall the construction of $`𝒯_n(\lambda )`$ given in Section 2.2. Let $`\omega `$ be an injective map from $`\{1,\mathrm{},m\}`$ to $`\{1,\mathrm{},n\}`$. Regarding $`\omega `$ as an $`m`$-tuple $`\omega =(\omega _1,\mathrm{},\omega _m)`$ we have
$$\epsilon _af(\omega )=\{\begin{array}{cc}0\hfill & \text{if }a\omega ,\hfill \\ f(\omega )\hfill & \text{if }a\omega .\hfill \end{array}$$
(4.58)
Therefore, the product $`(1\epsilon _{i_1})\mathrm{}(1\epsilon _{i_r})`$ is a projection to the subspace of functions $`f`$ such that the indices $`i_1,\mathrm{},i_r`$ belong to any $`\omega \mathrm{supp}f`$. This implies the first statement. The second follows from the obvious embedding $`H(\pi _\lambda )H(𝒯_n(\lambda ))`$ whose image consists of the functions supported by the maps $`\omega `$ such that $`\{\omega _1,\mathrm{},\omega _m\}=\{1,\mathrm{},m\}`$.
It was proved in (see also ) that the eigenvalue of $`c_n^{(r)}`$ in the irreducible representation $`\pi _\lambda `$ of $`S(n)`$ is a shifted symmetric function whose highest homogeneous component is the power sum symmetric function $`p_r`$.
###### Theorem 4.21
Let $`\mathrm{\Lambda }^{}`$ be identified with its image in $`\mathrm{Fun}`$. Then the mapping (4.57) is an isomorphism $`A_0\mathrm{\Lambda }^{}`$ of filtered algebras.
Proof. By Proposition 4.20 the images of the generators $`\mathrm{\Delta }^{(r)}A_0`$ with respect to the homomorphism (4.57) are shifted symmetric functions which are algebraically independent generators of the algebra $`\mathrm{\Lambda }^{}`$. The map obviously respects the filtrations.
Recall that by Propositions 3.10 and 3.11, elements of the center $`A_0`$ act in irreducible tame representations of $`S(\mathrm{})`$ by scalar operators. Hence, any such representation determines a homomorphism $`A_0`$.
###### Corollary 4.22
The center $`A_0`$ separates irreducible tame representations of $`S(\mathrm{})`$. That is, non-equivalent irreducible tame representations give rise to distinct homomorphisms $`A_0`$.
Proof. By Theorem 2.10, the irreducible tame representations are precisely the representations $`T_\lambda `$. Hence, our claim is equivalent to the fact that the map $`\mathrm{\Lambda }^{}\mathrm{Fun}`$ defined above is an embedding.
## 5 The structure of the algebra $`A_m`$, $`m>0`$
Here we generalize the results of Section 4 to the algebra $`A_m`$, where $`m=1,2,\mathrm{}`$. Throughout the section we assume $`0mn`$ and use the notation
$$_{mn}=\{m+1,\mathrm{},n\}.$$
(5.1)
For $`\gamma \mathrm{\Gamma }(n)`$, set
$$\begin{array}{cc}\hfill J_m(\gamma )& =\{i|i_{mn},\gamma _{ii}=0\},\hfill \\ \hfill \mathrm{deg}_m\gamma & =|J_m(\gamma )|.\hfill \end{array}$$
(5.2)
We shall call $`\mathrm{deg}_m\gamma `$ the $`m`$-degree of $`\gamma `$.
###### Proposition 5.1
For $`\gamma ,\delta \mathrm{\Gamma }(n)`$,
$$\mathrm{deg}_m\gamma \delta \mathrm{deg}_m\gamma +\mathrm{deg}_m\delta .$$
(5.3)
Proof. For any $`i_{mn}`$ we have
$$(\gamma \delta )_{ii}=0\gamma _{ij}\delta _{ji}=0\text{for all }j=1,\mathrm{},n\text{.}$$
(5.4)
In particular, $`(\gamma \delta )_{ii}=0`$ implies $`\gamma _{ii}\delta _{ii}=0`$, i.e.,
$$J_m(\gamma \delta )J_m(\gamma )J_m(\delta ),$$
(5.5)
and (5.3) follows.
###### Definition 5.2
Using the $`m`$-degree we define a new filtration in $`A(n)`$, called the $`m`$-filtration, by
$$A(m)=F_m^0(A(n))F_m^1(A(n))\mathrm{}F_m^{nm}(A(n))=A(n).$$
(5.6)
Here $`F_m^M(A(n))`$, the $`M`$-th term of the filtration, is formed by the elements $`aA(n)`$ which are linear combinations of the elements of $`\mathrm{\Gamma }(n)`$ of $`m`$-degree $`M`$. For any subspace $`S`$ of $`A(n)`$ we will use the symbol $`F_m^M(S)`$ to indicate the $`M`$-th term of the induced filtration.
By Proposition 5.1, the $`m`$-filtration is compatible with the algebra structure of $`A(n)`$, so the corresponding graded algebra exists. But contrary to the case $`m=0`$, this graded algebra is not commutative for $`m1`$ since it contains, as the 0-component, the non-commutative algebra $`A(m)`$.
Let $`D`$ be a multiplicative semigroup with unity 1. Consider the union $`D\{0\}`$, where 0 is an extra symbol, and adopt the convention that
$$d\mathrm{\hspace{0.17em}0}=0d=0,d+0=0+d=d\text{for any }dD.$$
(5.7)
###### Definition 5.3
(i) The semigroup $`S(m,D)`$ consists of the $`m\times m`$ matrices $`\alpha =[\alpha _{ij}]`$ with entries in $`D\{0\}`$ such that any row and column contains exactly one non-zero entry. The product is the matrix multiplication with the conventions (5.7).
(ii) The semigroup $`\mathrm{\Gamma }(n,D)`$ is defined as in (i) by allowing any row and column contain at most one non-zero entry.
Note that if $`D=\{1\}`$, then $`S(n,D)`$ and $`\mathrm{\Gamma }(n,D)`$ coincide with $`S(n)`$ and $`\mathrm{\Gamma }(n)`$, respectively. If $`D`$ is a group, then $`S(n,D)`$ is the wreath product of $`S(n)`$ and $`D`$.
We shall be assuming now that $`D`$ is the free abelian semigroup $`\{1,z,z^2,\mathrm{}\}`$ with unity 1 and one generator $`z`$. This semigroup is isomorphic to the additive semigroup $`_+`$. We denote the corresponding semigroups introduced in Definition 5.3 by $`S(m,_+)`$ and $`\mathrm{\Gamma }(m,_+)`$.
Set $`\mathrm{ord}z^k=k`$ for $`k=0,1,\mathrm{}`$, and for $`\alpha \mathrm{\Gamma }(m,_+)`$, set
$$\mathrm{ord}\alpha =\underset{i,j;\alpha _{ij}0}{}\mathrm{ord}\alpha _{ij}.$$
(5.8)
###### Definition 5.4
(i) Set
$$\mathrm{\Gamma }(m,n)=\{\sigma \mathrm{\Gamma }(n)|\mathrm{dom}\sigma \text{ and }\mathrm{range}\sigma \text{ contain }_{mn}\}.$$
(5.9)
(ii) Consider the linear span of $`\mathrm{\Gamma }(m,n)`$ and let $`Z_m(n)A(n)`$ denote the subspace in this span formed by the elements invariant under the conjugation by the elements of the group $`S_m(n)`$.
In particular, $`\mathrm{\Gamma }(0,n)=S(n)`$ and $`Z_0(n)=Z\left(S(n)\right)`$ is the center of $`[S(n)]`$. The role of $`\mathrm{\Gamma }(m,n)`$ and $`Z_m(n)`$ will be similar to that of $`S(n)`$ and $`Z\left(S(n)\right)`$ in Section 4. Note also that $`Z_m(n)`$ contains $`[S(n)]^{S_m(n)}`$, the centralizer of $`S_m(n)`$ in the group algebra $`[S(n)].`$
Now our purpose is to construct a convenient basis in $`Z_m(n)`$. To do this, we need to classify the $`S_m(n)`$-orbits in $`\mathrm{\Gamma }(m,n)`$ where the elements of $`S_m(n)`$ act by conjugations.
###### Proposition 5.5
There is a natural parameterization of the $`S_m(n)`$-orbits in $`\mathrm{\Gamma }(m,n)`$ by the couples $`(\alpha ,)`$, where $`\alpha \mathrm{\Gamma }(m,_+)`$ and $``$ is a partition such that
$$\mathrm{ord}\alpha +||=nm.$$
(5.10)
Proof. Fix an arbitrary element $`\sigma \mathrm{\Gamma }(m,n)`$ and assign to it an $`m\times m`$-matrix $`\alpha =\alpha (\sigma )`$ as follows. For $`i,j_{mn}`$ set
$`\alpha _{ij}=0`$ if $`j\mathrm{dom}\sigma `$, (5.11)
$`\alpha _{ij}=1`$ if $`j\mathrm{dom}\sigma `$ and $`\sigma (j)=i`$, (5.12)
$`\alpha _{ij}=z^k`$ $`\text{if }j\mathrm{dom}\sigma ,`$ (5.13)
and there exist $`k`$ points $`p_1,\mathrm{},p_k_{mn}`$ such that $`\sigma (j)=p_1`$, $`\sigma (p_1)=p_2`$, $`\mathrm{}`$, $`\sigma (p_{k1})=p_k`$, $`\sigma (p_k)=i`$. Thus, to any $`j\mathrm{dom}\sigma `$ with $`\sigma (j)_{mn}`$ we have assigned a subset $`\{p_1,\mathrm{},p_k\}_{mn}`$. It is clear that these subsets are pairwise disjoint. Let $`P=P(\sigma )`$ denote their union. Then $`\mathrm{ord}\alpha =|P|nm.`$ It is also clear that $`\alpha \mathrm{\Gamma }(m,_+)`$.
Further, let $`P^{}=P^{}(\sigma )`$ be the complement of $`P`$ in $`_{mn}`$. Then $`P^{}`$ is contained in the domain of $`\sigma `$, and $`P^{}`$ is $`\sigma `$-invariant. Therefore, the restriction of $`\sigma `$ to $`P^{}`$ defines a permutation of $`P^{}`$. Let $`=(\sigma )`$ be the partition of the number $`|P^{}|`$ which is defined by the lengths of the cycles of this permutation. Then the couple
$$(\alpha ,)=(\alpha (\sigma ),(\sigma ))$$
(5.14)
satisfies (5.10). It is clear that the couple (5.14) remains unchanged if $`\sigma `$ is replaced by $`s\sigma s^1`$ with $`sS_m(n)`$. Moreover, it is also clear that if the couples (5.14) corresponding to two elements of $`\mathrm{\Gamma }(m,n)`$ are the same, then these elements belong to the same orbit. Finally, any couple satisfying (5.10) can be obtained from an element of $`\mathrm{\Gamma }(m,n)`$.
Remark. A couple (5.14) corresponds to an element of $`S(n)\mathrm{\Gamma }(m,n)`$ if and only if $`\alpha S(m,_+)`$.
We shall now define analogs of the elements $`c_n^{}`$. First, for any subset $`Q_{mn}`$ and any partition $`=(M_1,\mathrm{},M_r)`$ such that $`||=|Q|`$ we set
$$c_Q^{}=(i_1,\mathrm{},i_{M_1})(j_1,\mathrm{},j_{M_2})\mathrm{}(k_1,\mathrm{},k_{M_r}),$$
(5.15)
where $`(i_1,\mathrm{},i_{M_1})`$ etc. are cyclic permutations of the corresponding indices and the summation is taken over all the orderings $`(i_1,\mathrm{},i_{M_1};j_1,\mathrm{},j_{M_2};\mathrm{};k_1,\mathrm{},k_{M_r})`$ of the elements of $`Q`$. We shall suppose that $`c_{\mathrm{}}^{\mathrm{}}=1`$.
Second, for any $`\alpha \mathrm{\Gamma }(m,_+)`$ and any subset $`P_{mn}`$ such that $`\mathrm{ord}\alpha =|P|`$, we set
$$\mathrm{\Gamma }(\alpha ,P)=\{\sigma \mathrm{\Gamma }(m,n)|\alpha (\sigma )=\alpha ,P(\sigma )=P,(\sigma )=(1^{nm|P|})\},$$
(5.16)
i.e., $`\sigma `$ has to fix all the points in $`_{mn}P`$.
###### Definition 5.6
For any couple $`(\alpha ,)`$, where $`\alpha \mathrm{\Gamma }(m,_+)`$ and $``$ is a partition such that $`\mathrm{ord}\alpha +||nm`$ we set
$$c_n^{\alpha ,}=\underset{P,Q}{}\underset{\sigma \mathrm{\Gamma }(\alpha ,P)}{}\sigma c_Q^{},$$
(5.17)
where $`P,Q`$ are disjoint subsets in $`_{mn}`$ such that
$$|P|=\mathrm{ord}\alpha ,|Q|=||.$$
(5.18)
###### Proposition 5.7
Each of the families
$$c_n^{\alpha ,}\text{with}\mathrm{ord}\alpha +||=nm,$$
(5.19)
and
$$c_n^{\alpha ,}\text{with}\mathrm{ord}\alpha +||nm,\text{and }\text{ has no part equal to }1,$$
(5.20)
forms a basis of $`Z_m(n)`$.
Proof. Note that $`c_n^{\alpha ,1\mathrm{}1}`$ is proportional to $`c_n^{\alpha ,}`$. Therefore, it suffices to consider the family (5.19). By Proposition 5.5 the elements $`c_n^{\alpha ,}`$ with $`\mathrm{ord}\alpha +||=nm`$ are proportional to characteristic functions of the $`S_m(n)`$-orbits in $`\mathrm{\Gamma }(m,n)`$.
Now we introduce analogs of the elements $`\mathrm{\Delta }_n^{}`$.
###### Definition 5.8
For any couple $`(\alpha ,)`$, where $`\alpha \mathrm{\Gamma }(m,_+)`$ and $``$ is a partition such that $`\mathrm{ord}\alpha +||nm`$ we set
$$\mathrm{\Delta }_n^{\alpha ,}=\underset{P,Q}{}\underset{\sigma \mathrm{\Gamma }(\alpha ,P)}{}\epsilon (P)\sigma c_Q^{}\epsilon (Q)\epsilon (P),$$
(5.21)
where $`\epsilon (I):=(1\epsilon _{i_1})\mathrm{}(1\epsilon _{i_k})`$ for $`I=\{i_1,\mathrm{},i_k\}`$. Here $`P,Q`$ are disjoint subsets in $`_{mn}`$ satisfying (5.18). We set $`\mathrm{\Delta }_n^{1,\mathrm{}}=1`$, where $`\mathrm{}`$ stands for the empty partition.
Note that (5.21) can be written in an equivalent form where the term $`\epsilon (Q)`$ takes the leftmost position; cf. (4.13) and (4.16).
###### Proposition 5.9
The elements $`\mathrm{\Delta }_n^{\alpha ,}`$ belong to the algebra $`A_m(n)`$.
Proof. The semigroup $`\mathrm{\Gamma }_m(n)`$ is generated by the subgroup $`S_m(n)`$ and the idempotents $`\epsilon _{m+1},\mathrm{},\epsilon _n`$. Therefore, it suffices to check that $`\mathrm{\Delta }_n^{\alpha ,}`$ is stable under the conjugation by the elements of $`S_m(n)`$ and commutes with the idempotents. The first claim is immediate from (5.21). The second claim is verified exactly as its counterpart for the elements $`\mathrm{\Delta }_n^{}`$; see the proof of Proposition 4.4.
The following is an analog of Proposition 4.5 and it is proved by the same argument.
###### Proposition 5.10
We have
$$\theta _n(\mathrm{\Delta }_n^{\alpha ,})=\mathrm{\Delta }_{n1}^{\alpha ,},$$
(5.22)
where we adopt the convention that
$$\mathrm{\Delta }_k^{\alpha ,}=0\text{if }\mathrm{ord}\alpha +||>km\text{.}$$
(5.23)
Our aim now is to prove an analog of Propositions 4.1 and 4.10; see Proposition 5.14 below. We need the following three lemmas.
###### Lemma 5.11
For $`m<n`$
$$I(n)A_m(n)=(1\epsilon _{m+1})\mathrm{}(1\epsilon _n)Z_m(n)(1\epsilon _{m+1})\mathrm{}(1\epsilon _n).$$
(5.24)
Proof. Suppose that $`xA(n)`$ can be written as
$$x=(1\epsilon _{m+1})\mathrm{}(1\epsilon _n)y(1\epsilon _{m+1})\mathrm{}(1\epsilon _n),$$
(5.25)
where $`yZ_m(n)`$. Then $`xA_m(n)`$ since $`x`$ is invariant under the conjugation by the elements of $`S_m(n)`$ and is annihilated when multiplied (from the left or from the right) by any idempotent $`\epsilon _{m+1},\mathrm{},\epsilon _n`$. Moreover, this also implies that $`xI(n)`$.
Conversely, suppose $`xI(n)A_m(n)`$. Then $`x\epsilon _n=\epsilon _nx=0`$. Using the invariance of $`x`$ under the conjugation by the elements of $`S_m(n)`$ we obtain $`x\epsilon _i=\epsilon _ix=0`$ for $`i=m+1,\mathrm{},n`$. Thus $`x`$ is invariant under the multiplication by $`(1\epsilon _{m+1})\mathrm{}(1\epsilon _n)`$ both from the left and from the right.
Further, we can write $`x=y+y^{}`$ where $`y`$ and $`y^{}`$ are spanned by elements of $`\mathrm{\Gamma }(m,n)`$ and $`\mathrm{\Gamma }(n)\mathrm{\Gamma }(m,n)`$, respectively. However,
$$(1\epsilon _{m+1})\mathrm{}(1\epsilon _n)y^{}(1\epsilon _{m+1})\mathrm{}(1\epsilon _n)=0$$
(5.26)
since for each element $`\gamma \mathrm{\Gamma }(n)\mathrm{\Gamma }(m,n)`$ there exists $`i>m`$ such that $`\gamma \epsilon _i=\gamma `$ or $`\epsilon _i\gamma =\gamma `$. This implies
$$x=(1\epsilon _{m+1})\mathrm{}(1\epsilon _n)y(1\epsilon _{m+1})\mathrm{}(1\epsilon _n).$$
(5.27)
Finally, averaging over the group $`S_m(n)`$ transforms $`y`$ into an element of $`Z_m(n)`$; cf. the proof of Lemma 4.7.
For $`\sigma \mathrm{\Gamma }(n)`$, set
$$Q(\sigma )=\{i_{mn}|\sigma _{ii}=1\}.$$
(5.28)
###### Lemma 5.12
The mapping
$$\sigma \gamma ,\gamma =\sigma \epsilon _{Q(\sigma )}=\epsilon _{Q(\sigma )}\sigma $$
(5.29)
defines a bijection of $`\mathrm{\Gamma }(m,n)`$ onto the set of all $`\gamma \mathrm{\Gamma }(n)`$ satisfying the conditions
$`\mathrm{dom}\gamma _{mn}`$ $`=\mathrm{range}\gamma _{mn},`$ (5.30)
$`\mathrm{deg}_m\gamma `$ $`=nm.`$ (5.31)
Proof. The effect of the multiplication of $`\sigma `$ by $`\epsilon _{Q(\sigma )}`$ from the left or from the right consists of replacing all the diagonal entries $`\sigma _{ii}=1`$ with $`i>m`$ by zeros. Therefore $`\gamma `$ satisfies (5.30). Relation (5.31) follows from this observation and the fact that both $`\mathrm{dom}\sigma `$ and $`\mathrm{range}\sigma `$ contain $`_{mn}`$.
Conversely, let $`\gamma \mathrm{\Gamma }(n)`$ satisfy (5.30) and (5.31). Note that (5.30) can be reformulated as follows: for any $`i=m+1,\mathrm{},n`$ the $`i`$-th row and the $`i`$-th column are zero or non zero at the same time, whereas (5.31) means that all the diagonal entries $`\gamma _{ii}`$ with $`i>m+1`$ vanish. Now, let $`\sigma `$ be defined by
$`\sigma _{ij}`$ $`=\gamma _{ij}`$ $`\text{if either }ij\text{ or }\mathrm{min}\{i,j\}m,`$
$`\sigma _{ii}`$ $`=1`$ if $`i_{mn}`$
and the $`i`$-th row (or the $`i`$-th column) of $`\gamma `$ is zero. Then it is easy to see that $`\sigma \mathrm{\Gamma }(m,n)`$ and that $`\gamma `$ is the image of $`\sigma `$ under the mapping (5.29).
###### Lemma 5.13
For $`m<n`$ the restriction of the projection $`\theta _n:A_m(n)A_m(n1)`$ to the subspace $`F_m^{nm1}(A_m(n))`$ is injective.
Proof. Suppose that $`xA_m(n)`$ and $`\theta _n(x)=0`$. We will show that then $`\mathrm{deg}_mx=nm`$ unless $`x=0`$.
By Lemma 5.11, $`x`$ can be written as a linear combination of the elements of type
$$\begin{array}{c}(1\epsilon _{m+1})\mathrm{}(1\epsilon _n)\sigma (1\epsilon _{m+1})\mathrm{}(1\epsilon _n)\hfill \\ \hfill =\underset{R,SN_{mn}}{}(1)^{|R|+|S|}\epsilon _R\sigma \epsilon _S,\sigma \mathrm{\Gamma }(m,n).\end{array}$$
(5.32)
Let us divide the terms in the sum (5.32) into two groups depending on whether $`RS`$ contains $`Q(\sigma )`$ or not. Then the terms of the first group are of $`m`$-degree $`nm`$ whereas those of the second group are of $`m`$-degree $`<nm`$. So, it suffices to prove that the elements
$$\underset{RSQ(\sigma )}{}(1)^{|R|+|S|}\epsilon _R\sigma \epsilon _S,\sigma \mathrm{\Gamma }(m,n),$$
(5.33)
are linearly independent. Note that, in the case $`RS=Q(\sigma )`$,
$$\epsilon _R\sigma \epsilon _S=\sigma \epsilon _{Q(\sigma )}\text{and}\mathrm{rank}\sigma \epsilon _{Q(\sigma )}=\mathrm{rank}\sigma |Q(\sigma )|,$$
(5.34)
whereas, in the case $`RS`$ strictly contains $`Q(\sigma )`$,
$$\mathrm{rank}\sigma \epsilon _{Q(\sigma )}<\mathrm{rank}\sigma |Q(\sigma )|.$$
(5.35)
Therefore, we now need to show that for any fixed $`k`$ the elements
$$\underset{RS=Q(\sigma )}{}(1)^{|R|+|S|}\sigma \epsilon _{Q(\sigma )},$$
(5.36)
where $`\sigma `$ runs over the subset of the elements in $`\mathrm{\Gamma }(m,n)`$ with $`\mathrm{rank}\sigma |Q(\sigma )|=k`$, are linearly independent.
Lemma 5.12 implies that the elements $`\sigma \epsilon _{Q(\sigma )}\mathrm{\Gamma }(n)`$ are pairwise distinct. Hence it remains to prove that all the coefficients in (5.36) are non-vanishing. This is implied by the following general fact: if $`Q`$ is an arbitrary finite set, then
$$\underset{R,SQ,RS=Q}{}(1)^{|R|+|S|}0.$$
(5.37)
We will prove that the sum in (5.37) equals $`(1)^q`$ where $`q=|Q|`$. Indeed, for any $`r=0,1,\mathrm{},q`$, there are $`\left({\displaystyle \genfrac{}{}{0pt}{}{n}{r}}\right)`$ subsets $`RQ`$ with $`|R|=r`$. Given $`R`$, for any $`t=0,1,\mathrm{},r`$, there are $`\left({\displaystyle \genfrac{}{}{0pt}{}{r}{t}}\right)`$ subsets $`SQ`$ such that $`RS=Q`$ and $`|RS|=t`$. Since
$$|R|+|S|=r+t+(qr)=q+t,$$
(5.38)
the sum in (5.37) equals
$$(1)^q\underset{r=0}{\overset{q}{}}\left(\genfrac{}{}{0pt}{}{n}{r}\right)\underset{t=0}{\overset{r}{}}(1)^t\left(\genfrac{}{}{0pt}{}{r}{t}\right).$$
If $`r=0`$ then the interior sum is equal to $`1`$, otherwise it is zero. Therefore the entire sum is $`(1)^q`$.
###### Proposition 5.14
The elements $`\mathrm{\Delta }_n^{\alpha ,}`$ with
$$\mathrm{ord}\alpha +||nm$$
(5.39)
form a basis of $`A_m(n)`$. Moreover, for any $`M`$ with $`0Mnm`$ the elements $`\mathrm{\Delta }_n^{\alpha ,}`$ satisfying
$$\mathrm{ord}\alpha +||M$$
(5.40)
form a basis of $`F_m^M(A_m(n))`$.
Proof. It suffices to prove the second claim. We use induction on $`n`$ and follow the argument of the proof of Proposition 4.10. The claim is obviously true for $`n=m`$. Assume that $`nm+1`$ and $`Mnm1`$. Lemma 5.13 implies that the elements $`\mathrm{\Delta }_n^{\alpha ,}`$ with $`\mathrm{ord}\alpha +||M`$ form a basis of $`F_m^M(A_m(n))`$.
To show that the elements $`\mathrm{\Delta }_n^{\alpha ,}`$ with $`\mathrm{ord}\alpha +||=nm`$ form a basis of $`I(n)A_m(n)`$ note that
$$\mathrm{\Delta }_n^{\alpha ,}=(1\epsilon _{m+1})\mathrm{}(1\epsilon _n)c_n^{\alpha ,}(1\epsilon _{m+1})\mathrm{}(1\epsilon _n);$$
(5.41)
see (5.21). Now the claim follows from Proposition 5.7 and the fact that the elements $`c_n^{\alpha ,}`$, being multiplied by $`(1\epsilon _{m+1})\mathrm{}(1\epsilon _n)`$, remain linearly independent; cf. (4.38).
Using Proposition 5.10 we can introduce the elements $`\mathrm{\Delta }^{\alpha ,}A_m`$ as sequences $`\mathrm{\Delta }^{\alpha ,}=(\mathrm{\Delta }_n^{\alpha ,}|nm)`$.
Remark. We can regard $`\mathrm{\Delta }^{\alpha ,}`$ as a formal series given by (5.21) where the sum is taken over all disjoint subsets $`P`$ and $`Q`$ in $`\{m+1,m+2,\mathrm{}\}`$ satisfying (5.18).
###### Theorem 5.15
The elements $`\mathrm{\Delta }^{\alpha ,}`$ with $`\alpha \mathrm{\Gamma }(m,_+)`$ and $``$ form a basis of the algebra $`A_m`$. Moreover, for any $`M0`$, the elements $`\mathrm{\Delta }^{\alpha ,}`$ with $`\mathrm{ord}\alpha +||M`$ form a basis of the $`M`$-th subspace $`F_m^M(A_m)`$ in $`A_m`$.
Proof. The first claim follows from the second one. The second claim follows from Proposition 5.14 and the definition of $`F_m^M(A_m)`$ as the projective limit of the spaces $`F_m^M(A_m(n))`$.
###### Corollary 5.16
For $`n>M`$, the mapping
$$\theta _n:F_m^M(A_m(n))F_m^M(A_m(n1))$$
(5.42)
is an isomorphism of vector spaces and so is the mapping
$$\theta ^{(n)}:F_m^M(A_m)F_m^M(A_m(n)),nM.$$
(5.43)
In particular, $`dimF_m^M(A_m)<\mathrm{}`$.
For each $`k=1,\mathrm{},m`$ consider the following elements of $`A_m(n)`$
$$u_{k|n}=\underset{i=k+1}{\overset{n}{}}(ki)(1\epsilon _k)(1\epsilon _i)=\underset{i=k+1}{\overset{n}{}}(1\epsilon _i)(ki)(1\epsilon _i).$$
(5.44)
The image of $`u_{k|n}`$ under the retraction homomorphism (3.9) is the Jucys–Murphy element for $`S(n)`$; cf. , . We obviously have $`\theta _n(u_{k|n})=u_{k|n1}`$ and so, for each $`k`$ the element $`u_kA_m`$ can be defined as the sequence $`u_k=(u_{k|n}|nm)`$. Recall that the algebra $`A(m)`$ is naturally embedded in $`A_m`$; see Proposition 3.7.
###### Proposition 5.17
The following relations hold in the algebra $`A_m`$:
$`s_ku_k`$ $`=u_{k+1}s_k+(1\epsilon _k)(1\epsilon _{k+1}),`$ $`s_ku_l`$ $`=u_ls_k,lk,k+1;`$ (5.45)
$`u_ku_l`$ $`=u_lu_k,\epsilon _ku_k=u_k\epsilon _k=0,`$ $`\epsilon _iu_k`$ $`=u_k\epsilon _i,ik;`$ (5.46)
where $`s_k=(k,k+1)`$.
Proof. For $`n>m`$ we have $`u_{1|n}=\mathrm{\Delta }_n^{(2)}\mathrm{\Delta }_{n1}^{(2)}`$, where $`\mathrm{\Delta }_{n1}^{(2)}`$ is the element of the center of $`[\mathrm{\Gamma }_1(n)]`$ given by (4.13) with the sum taken over the indices from $`\{2,\mathrm{},n\}`$. Now an easy induction proves that the elements $`u_{1|n},\mathrm{},u_{m|n}`$ pairwise commute, and so do the elements $`u_1,\mathrm{},u_m`$. The remaining relations easily follow from (5.44) and the relations in the algebra $`A(n)`$.
We shall denote by $`\stackrel{~}{}_m`$ the subalgebra of $`A_m`$ generated by $`A(m)`$ and the elements $`u_1,\mathrm{},u_m`$. The following is our main result. The theorem describes the structure of the algebra $`A_m`$.
###### Theorem 5.18
We have an algebra isomorphism
$$A_mA_0\stackrel{~}{}_m.$$
(5.47)
Moreover, the algebra $`\stackrel{~}{}_m`$ is isomorphic to an abstract algebra with generators
$`s_1,\mathrm{},s_{m1}`$, $`\epsilon _1,\mathrm{},\epsilon _m`$, $`u_1,\mathrm{},u_m`$ and the defining relations given by (3.2)–(3.4) and (5.45)–(5.46).
Recall that by Theorem 4.21 $`A_0`$ is isomorphic to the algebra of shifted symmetric functions $`\mathrm{\Lambda }^{}`$.
Proof. For any $`\alpha \mathrm{\Gamma }(m,_+)`$ and $``$ such that $`\mathrm{ord}\alpha +||nm`$ we have the equality in the algebra $`A_m(n)`$,
$$\mathrm{\Delta }_n^{\alpha ,\mathrm{}}\mathrm{\Delta }_n^{1,}=\mathrm{\Delta }_n^{\alpha ,}+\text{ lower }m\text{-degree terms},$$
(5.48)
where $`\mathrm{}`$ stands for the empty partition while $`1\mathrm{\Gamma }(m,_+)`$ is the $`m\times m`$ identity matrix; cf. the proofs of Proposition 4.2 and Corollary 4.11. On the other hand, we have
$$\mathrm{deg}_m(\mathrm{\Delta }_n^{}\mathrm{\Delta }_n^{1,})<||;$$
(5.49)
see (4.13) for the definition of $`\mathrm{\Delta }_n^{}`$. Now (5.48) and Proposition 5.14 imply that the elements $`\mathrm{\Delta }_n^{\alpha ,\mathrm{}}\mathrm{\Delta }_n^{}`$ with $`\mathrm{ord}\alpha +||nm`$ form a basis of $`A_m(n)`$. Hence the elements $`\mathrm{\Delta }^{\alpha ,\mathrm{}}\mathrm{\Delta }^{}`$ with $`\alpha \mathrm{\Gamma }(m,_+)`$ and $``$ form a basis of the algebra $`A_m`$. In other words, $`A_m`$ is a free $`A_0`$-module with the basis $`\{\mathrm{\Delta }^{\alpha ,\mathrm{}}|\alpha \mathrm{\Gamma }(m,_+)\}`$.
Further, if $`\alpha \mathrm{\Gamma }(m,_+)`$ has zero rows $`i_1,\mathrm{},i_r`$ then
$$\mathrm{\Delta }_n^{\alpha ,\mathrm{}}=\epsilon _{i_1}\mathrm{}\epsilon _{i_r}\mathrm{\Delta }_n^{\alpha ^{},\mathrm{}}$$
(5.50)
for some element $`\alpha ^{}S(m,_+)`$. Observe now that every element $`\alpha ^{}S(m,_+)`$ can be written as a product of the form
$$\alpha ^{}=\sigma \alpha _1^{k_1}\mathrm{}\alpha _m^{k_m},k_i0,$$
(5.51)
where $`\sigma S(m)`$ and $`\alpha _i\mathrm{\Gamma }(m,_+)`$ is the diagonal matrix whose $`ii`$-th entry is $`z`$ and all other diagonal entries are equal to $`1`$. This implies that modulo lower $`m`$-degree terms, the element $`\mathrm{\Delta }_n^{\alpha ^{},\mathrm{}}`$ coincides with the product
$$\mathrm{\Delta }_n^{\sigma ,\mathrm{}}(\mathrm{\Delta }_n^{\alpha _1,\mathrm{}})^{k_1}\mathrm{}(\mathrm{\Delta }_n^{\alpha _m,\mathrm{}})^{k_m};$$
(5.52)
cf. the proof of (4.44). The claim remains valid if we replace each $`\mathrm{\Delta }_n^{\alpha _k,\mathrm{}}`$ with the element $`u_{k|n}`$. Indeed, this follows from the equality
$$\mathrm{\Delta }_n^{\alpha _k,\mathrm{}}=u_{k|n}+\text{ elements of }m\text{-degree zero}.$$
(5.53)
Note also that the element $`\mathrm{\Delta }_n^{\sigma ,\mathrm{}}`$ can be identified with $`\sigma `$. Thus, modulo lower $`m`$-degree terms, the element (5.52) coincides with the product $`\sigma u_{1|n}^{k_1}\mathrm{}u_{m|n}^{k_m}`$. Using an obvious induction on the $`m`$-degree we may conclude that the $`A_0`$-module $`A_m`$ is generated by the subspace $`\stackrel{~}{}_m`$.
To prove that $`\stackrel{~}{}_m`$ generates the $`A_0`$-module $`A_m`$ freely, we check that for any $`M>0`$ the dimension of the subspace $`F_m^M(\stackrel{~}{}_m)`$ is less or equal to the number of elements $`\alpha \mathrm{\Gamma }(m,_+)`$ with $`\mathrm{ord}\alpha M`$. Indeed, by Proposition 5.17 the subalgebra $`\stackrel{~}{}_m`$ is spanned by the elements of the form $`\gamma u_1^{k_1}\mathrm{}u_m^{k_m}`$ with $`\gamma \mathrm{\Gamma }(m)`$. The relation $`\epsilon _ku_k=0`$ ensures that such a product is zero unless $`k_j=0`$ for each zero column $`j`$ in $`\gamma `$. To each of the nonzero products associate the element $`\alpha \mathrm{\Gamma }(m,_+)`$ which has the $`ij`$-entry $`z^{k_j}`$ where the $`j`$-th column of $`\gamma `$ is nonzero with $`\gamma _{ij}=1`$. This shows that the cardinality of a basis of $`F_m^M(\stackrel{~}{}_m)`$ can be at most the number of elements $`\alpha \mathrm{\Gamma }(m,_+)`$ with $`\mathrm{ord}\alpha M`$, proving (5.47).
To prove the second claim of the theorem note that by Proposition 5.17 there is an algebra epimorphism from the abstract algebra in question to $`\stackrel{~}{}_m`$. The above argument implies that the nonzero products $`\gamma u_1^{k_1}\mathrm{}u_m^{k_m}`$ with $`\gamma \mathrm{\Gamma }(m)`$ form a basis of $`\stackrel{~}{}_m`$.
###### Corollary 5.19
The mapping
$$s_ks_k,\epsilon _k\epsilon _k,u_ku_{k|n}$$
(5.54)
defines an algebra homomorphism $`\psi :\stackrel{~}{}_mA_m(n)`$. The algebra $`A_m(n)`$ is generated by $`A_0(n)`$ and the image of $`\psi `$.
The degenerate affine Hecke algebra $`_m`$ (see , ) is defined to be generated by elements $`s_1,\mathrm{},s_{m1}`$ and $`u_1,\mathrm{},u_m`$ with the defining relations (3.2) and
$`s_ku_k`$ $`=u_{k+1}s_k+1,`$ $`s_ku_l`$ $`=u_ls_k,lk,k+1;`$ (5.55)
$`u_ku_l`$ $`=u_lu_k.`$ (5.56)
As a linear space, $`_m`$ is isomorphic to the tensor product $`[S(m)][u_1,\mathrm{},u_m]`$. The following corollary is implied by Theorem 5.18 and provides an analog of the retraction homomorphism (3.9).
###### Corollary 5.20
The mapping
$$s_ks_k,u_ku_k,\epsilon _k0$$
(5.57)
defines an algebra epimorphism $`\stackrel{~}{}_m_m`$.
It can be seen from the proof of Theorem 5.18 that the retraction homomorphisms (3.9) and (5.57) “respect” the homomorphism $`\psi :\stackrel{~}{}_mA_m(n)`$ defined in Corollary 5.19. More precisely, the following result takes place. It was announced in \[29, Theorem 11\], and a proof was given in . We denote by $`B_m(n)`$ the centralizer of $`S_m(n)`$ in the group algebra $`[S(n)]`$; see Introduction.
###### Corollary 5.21
The mapping
$$s_ks_k,u_k\underset{i=k+1}{\overset{n}{}}(ki)$$
(5.58)
defines an algebra homomorphism $`\phi :_mB_m(n)`$. The algebra $`B_m(n)`$ is generated by $`B_0(n)`$ and the image of $`\phi `$. |
warning/0002/astro-ph0002469.html | ar5iv | text | # Revealing the Photodissociation Region: Hubble Space Telescope/NICMOS Imaging of NGC 7027
## 1 Introduction
Planetary nebulae have long been thought of as ionized remnants of circumstellar material ejected by stars on the asymptotic giant branch (AGB). That somewhat limited view has changed dramatically, as the neutral component of the remnant circumstellar envelope has been shown to remain observable very late into the lifetime of many planetary nebulae (PNe; e.g., Dinerstein, Sneden, & Uglum (1995); Huggins et al. (1996); Hora, Latter, & Deutsch (1999); and references therein). It is now known that many PNe are characterized by emission from material that is contained in rapidly evolving photodissocation regions (PDRs; Tielens & Hollenbach (1985); Sternberg & Dalgarno (1989); Latter, Walker, & Maloney (1993); Natta & Hollenbach (1998)). As the central star and nebula evolve, UV emission from the star increases rapidly. A photodissociation front moves through the gas and slowly turns the mostly molecular nebula to predominantly atomic. At this point, the nebula radiates mainly in the infrared. When the central star becomes hotter than $`T_{}30,000`$ K, the circumstellar gas quickly becomes ionized and the planetary nebula shines brightly in visible light. The chemical properties of gas in PNe and proto-PNe (PPN) are like those of interstellar PDRs (Latter et al. (1992); Latter, Walker, & Maloney (1993); Tielens (1993); Hora & Latter (1994)), which are well modeled (see Hollenbach & Tielens (1997)). The time over which molecular material can be present in PNe is typically predicted to be a rather brief period in the lifetime of PNe (e.g., Tielens (1993); Natta & Hollenbach (1998)). But, such timescales predicted from current chemical models are limited by observational data and lack of a clear understanding of the morphology and density structure of these objects.
NGC 7027 is perhaps the best studied planetary nebula, in part because of its proximity (about 900 pc; e.g. Masson (1989)) and high surface brightness at all wavelengths. It has a very rich atomic and molecular spectrum, making it ripe for study at all wavelengths, and it is rich in physical and chemical information (see, e.g., Graham et al. 1993a ; Cox et al. (1997), and references therein). Because of its compact size as viewed from Earth ($`15`$″), the morphology, particularly that for the H<sub>2</sub> emission, of this object has been suggested, but has not been clearly revealed (see, e.g., Graham et al. 1993a ; Graham et al. 1993b ; Kastner et al. (1994); Latter et al. (1995)). The elliptical ionized core lies at the center of an extended molecular envelope (Bieging, Wilner, & Thronson (1991)). At the interface between the cold molecular envelope and the hot ionized region is an apparent “quadrupolar” region of strong near–infrared (IR) rotational-vibrational molecular hydrogen emission (Graham et al. 1993a ) that has been shown to be excited by absorption of UV photons in the photodissociation region (Hora, Latter, & Deutsch (1999)).
A detailed understanding of the morphology of planetary nebulae is important. It has been demonstrated that there is a strong correlation between the presence of molecular emission from PNe and the observed morphology of the PNe, such that objects that contain large amounts of molecular material are bipolar or butterfly nebulae (Hora, Latter, & Deutsch (1999); Kastner et al. (1996); Huggins et al. (1996); Huggins & Healy (1989); Zuckerman & Gatley (1988)). Progenitor mass also correlates with morphological type, such that higher mass stars appear to produce bipolar or butterfly nebulae (see Corradi & Schwarz (1994)). This correlation suggests that the higher mass AGB progenitors, which probably have the highest mass loss rates, produce dense, long lived molecular envelopes (Hora, Latter, & Deutsch (1999)).
The evolution of PDR and PNe, and the correlation between morphology with molecular content are not fully understood. Because NGC 7027 is in a rapid and key moment in evolution (the transition from neutral, predominantly molecular envelope to an ionized one), it is worthy of serious study into its chemical and physical properties and its detailed morphology. The Hubble Space Telescope and NICMOS provided the high spatial resolution, dynamic range, and stable background needed to examine the near-IR nature of NGC 7027. In this paper, we present narrowband images in the $`v=10`$ S(1) ($`\lambda =2.121`$ µm) line of molecular hydrogen and other filters that trace the ionized core and the nearby neutral region. These images show with unprecedented clarity the true structure of NGC 7027. In Sections 2 and 3, we discuss the data, the apparent morphology, and the presence of a jet (or jets) in the object. In §4 we present a 3-dimensional model for the overall structure, including the photodissociation region. Section 5 presents a discussion of the excitation of the nebula and the central star properties. We then consider our results in light of previous work, and discuss the evolution of the PDR in NGC 7027.
## 2 Observations and Data Reduction
We have imaged NGC 7027 with the Hubble Space Telescope (HST) and the Near-Infrared Camera Multiobject Spectrometer (NICMOS) at 7 near-infrared wavelengths between 1.1 µm and 2.15 µm (Figs. 1 and 2). The 1.10 µm broadband (F110W) observations were made with Camera 1 (pixel scale = 0.043″ pixel<sup>-1</sup>, field of view or FOV = 11″$`\times `$11″). All of the other observations were made with Camera 2 (pixel scale = 0.075″ pixel<sup>-1</sup>, FOV = 19.2″$`\times `$19.2″). The observations were made in one of the multiaccum detector readout modes (usually multiaccum 256; see the NICMOS Instrument Handbook) to obtain a very high dynamic range. A variety of dither patterns were utilized to sample the array uniformly; in the case of the F110W observations dithering/mosaicing was required to obtain a complete image of the source. Chop (or blank sky) frames for background subtraction were also obtained for all filters except F110W and F160W (1.60 µm broadband). From ground-based observations it was known that the full extent of the H<sub>2</sub> emission in NGC 7027 is slightly larger than the NICMOS Camera 2 field of view. A specific orientation of the Observatory was used to put the largest known extent of the emission along a diagonal of the array. Dither patterns were used that would ensure complete coverage of the emission. The method worked very well, and no emission was missed, although it does appear that very low level scattered light does extend off the array in some of the bands observed. A summary of the observations is presented in Table 1, and the data are presented in Figs. 1 and 2.
The data were originally reduced and calibrated by the standard Space Telescope Science Institute (STScI) NICMOS pipeline, but we subsequently re-reduced all of the data using better characterized calibration files (flat-fields, dark frames, and bad pixel masks) as well as a revised version of the STSDAS calibration task calnicA. After reduction of individual frames, the multiple images at each wavelength were mosaiced together using the STSDAS calnicB task in IRAF<sup>3</sup><sup>3</sup>3The Image Reduction and Analysis Facility is written and supported by the National Optical Astronomy Observatories (NOAO) in Tucson, AZ. The reduced images were flux calibrated using NICMOS photometric calibration tables provided by M. Rieke (private communication). The HST Wide Field/Planetary Camera (WFPC2) data were acquired from the STScI HST data archive.
##### Isolating the H<sub>2</sub> emission –
An examination of the near-IR spectrum of NGC 7027 (Hora, Latter, & Deutsch (1999)), and the F212N and F215N filter transmission curves shows that both filters suffer from line contamination. In addition to the H<sub>2</sub> $`v=10`$ S(1) line at 2.121 µm, the F212N filter also transmits the 2.113 µm He I line at $``$ 83% of peak transmission. The F215N filter, used for sampling continuum emission adjacent to the H<sub>2</sub> line, includes within the filter bandpass the Br$`\gamma `$ line at 2.167 µm at $``$ 2 – 4% transmission. We made use of all available data, theoretical modeling, and empirical fitting to provide the cleanest possible subtraction of the continuum and contaminating lines in the H<sub>2</sub> image. A simple subtraction of the flux calibrated continuum image from the line image not only isolates the H<sub>2</sub> emission but also shows a dark (negative) inner ring-shaped region that resembles the ionized nebula; this inner region is an artifact of over–subtraction of the He I line contribution from the H<sub>2</sub> filter data. Typically, a bright ring of He I emission is seen in ground–based narrowband H<sub>2</sub> images of NGC 7027 (see, e.g., Graham et al. 1993a ; Latter et al. (1995); Kastner et al. (1994); Kastner et al. (1996)), because the continuum filters most often used do not include Br$`\gamma `$ emission in the bandpass.
Since uncertainties associated with filter transmission are largest along the edges of the bandpass (where transmission changes rapidly with wavelength), we first made an estimate of the Br$`\gamma `$ transmission. Using the F190N image as “true” continuum for the F215N image, we compared the observed Br$`\gamma `$ flux with that obtained by Hora, Latter, & Deutsch (1999) using ground–based spectroscopy. We measured our fluxes at both of the spatial locations at which their slits were placed (labelled N and NW in their paper) and found that we detect approximately 7 – 8% of their Br$`\gamma `$ fluxes at each point. Since the Br$`\gamma `$ to He I ratio is $``$ 15 (Hora, Latter, & Deutsch (1999)) and assuming that the He I line is transmitted at 83% of peak, this implies that we need to scale the F215N image by $``$ 0.8 to cancel out the He I emission in the (F212N – F215N) difference image. This process of eliminating line contamination assumes that the the He I to Br$`\gamma `$ ratio is fairly constant over the nebula. Photoionization modeling using CLOUDY (Ferland (1997)) shows that S<sub>HeI</sub>/S<sub>Brγ</sub> ratio does not change by more than 20% over the region 1.5$`\times `$10<sup>17</sup> cm to the outer edge, 3$`\times `$10<sup>17</sup> cm ($``$ 10″– 20″). Thus, we can scale the Br$`\gamma `$ emission (i.e., the F215N image) to cancel out the He I emission in F212N – F215N difference image to a very high degree. Given the uncertainties involved in this procedure, we used scale factors in the range 0.7 – 1 to determine the best H<sub>2</sub> difference image empirically. We found that a value of 0.9, when applied to the F215N image, provided the cleanest subtraction of the continuum and He I from the F212N image. This scaling factor is consistent with $``$ 6% Br$`\gamma `$ transmission in the F215N filter. This procedure has resulted in the nearly pure molecular hydrogen emission image of NGC 7027 shown in Figure 3.
The continuum subtracted H<sub>2</sub> image has an integrated line brightness of 3.8$`\times `$10<sup>-12</sup> erg s<sup>-1</sup> cm<sup>-2</sup> and an average surface brightness of 8$`\times `$10<sup>-4</sup> erg s<sup>-1</sup> cm<sup>-2</sup> ster<sup>-1</sup>. Using a similar process of aligning images before subtraction we find an integrated P$`\alpha `$ line flux (the F187N line filter and F190N continuum filter; see Fig. 1) of $`5.5\times `$10<sup>-10</sup> erg s<sup>-1</sup> cm<sup>-2</sup>(Fig. 1).
## 3 General Results
### 3.1 Observed Nebular Properties
Most of the images (narrow- and broad-band) are remarkably similar (see Figs. 1 and 2). The emission is dominated by a bright, patchy elliptical ring (major axis along PA = 150°) which is the ionized nebula; the overall shape and extent of this nebula is similar to that previously observed in visible light, near-IR, and radio continuum images (e.g., Roelfsema et al. (1991); Woodward et al. (1992); Graham et al. 1993a ). Compared to the visible light images of HST/WFPC2 that show a “filled-in” ellipsoid, the near-IR morphology (where optical depths are lower) clearly shows an ellipsoidal shell. The high spatial resolution of these new HST/NICMOS images clearly resolves the ellipsoidal shell into apparent clumps (or “clump-like” variations in brightness), regions of varying surface brightness, and dark dust lanes. The brightest region appears to be a clump located about 4″ NW of the center of the nebula (marked by the location of the visible central star). The central star stands out clearly in the continuum images. The similarity between the P$`\alpha `$ and continuum images suggest that the broad-band filters trace free-free or free-bound emission from hot gas, rather than dust-scattered starlight, which is apparent as extended emission in the WFPC2 data. The F212N filter clearly shows a second, spatially distinct component that is not seen in the other filters: wisps of ro-vibrationally excited H<sub>2</sub> emission that appears to surround the ionized nebula. Furthermore, there is also clear evidence of a disruption of the gas along a NW-SE (PA $``$ 130° measured E of N) axis; this might be caused by wind interactions with an off-axis bipolar jet (see §3.2).
The continuum–subtracted H<sub>2</sub> image (Fig. 3) clearly shows two bright intersecting rings, which we interpret as brightened rims of an inclined biconical structure (§4.1). This structure encloses the roughly elliptical ionized region. The structures seen in molecular hydrogen emission and in tracers of the ionized core appear different and spatially distinct.
### 3.2 A Possible Collimated, Off-axis Jet
The H<sub>2</sub> emission (Figs. 1 – 3) shows bubble-like extensions at the outer edge of the bulk emission region. This apparent disturbance falls along an axis $``$ 20°– 40° (PA = 130° NW to SE) to the west of the polar axis, extending through the central star to the opposite side of the nebula (PA = 130° NW to SE). There are hints of these structures in ground-based near–IR data (Graham et al. 1993a ; Kastner et al. (1994); Latter et al. (1995)) as well as in radio continuum images of the ionized region (Roelfsema et al. (1991)). Because of the clarity of these new data, we can now suggest a cause for this structure. The wisps of H<sub>2</sub> emission revealed clearly in these HST data indicate that this region has been disturbed by something other than dust–photon interactions or a uniform, radially directed, outflowing wind. We suggest that there is present in NGC 7027 a highly collimated, “off-axis” jet. The disturbance is not only seen in the molecular region, but is clearly evident in the ionized core along precisely the same position angle. The structure extends into the region of scattered emission seen in the WFPC2 images as well (see also Roelfsema et al. (1991)). The sharpness of the extended WFPC2 emission argues against something like turbulent instabilities. There is no kinematical evidence for the jet, nor is there any extended collimated emission in the HST images; this leads us to believe that the jet may have disrupted the nebula in the past and has now died down. At a much lower level, there appears to be similar evidence for another jet, or the same jet precessing, at a polar angle of $``$ 10° to the east of the north pole. Along an axis in this direction there is an additional region of “disturbed” H<sub>2</sub> emission most clearly seen on the southern side of the nebula. The ionized core has apparent breaks or disruptions along this axis on the northern side, and extensions can be seen in the WFPC2 data as well (Figs. 2 and 3). The locations of these features are indicated by lines 1 and 2 in Fig. 3.
We speculate that the disturbed morphology of NGC 7027 might be a result of one or more BRETs (bipolar rotating episodic jets; e.g., Lopez et al. (1998)). In this regard we note that the two axes that we have identified simply show the most pronounced structure; however much of the other billowy structure and wisps (especially prominent in the H<sub>2</sub> data) is likely caused by the same, or similar processes. Evidence for precessing, collimated jets that are on, or away from the primary PN axis are also seen in a number of other PNe. The most striking example might be NGC 6543 (the “Cat’s Eye Nebula;” Lame, Harrington & Borkowski (1997)). NGC 7027 appears to be in a very early stage of forming a jet, or perhaps the jet is a transient phenomenon. The fact that we do not see the jet itself suggests that either we have not found the right diagnostic line or wavelength to probe the jet kinematically, or the jet has now shut off.
## 4 A 3-dimensional Spatial and Kinematical Model
We have modeled the spatial and kinematical structure of NGC 7027 using a geometrical modeling code described by Dayal et al. (1999). The code uses cartesian coordinates (for adaptability to different geometries), and the nebula is divided into 10<sup>7</sup> 0$`{}_{}{}^{}._{}^{}`$075-sized cells. In these models the nebula is assumed to have azimuthal symmetry with the star at the center. Once the shape and dimensions of a source have been set, the density, temperature, or velocity fields can be specified independently as functions of the radial distance from the center and/or latitude. The expansion of the nebula is assumed to be purely radial, directed away from the central star. The tilt angle of the nebula (w.r.t. the plane of the sky) can be varied independently for each model. Assuming that the emission is optically thin, the code numerically integrates the emission per cell along each line of sight, to produce surface brightness and position–velocity (P–V) maps. The surface–brightness maps are convolved with a gaussian beam (FWHM = 0$`{}_{}{}^{}._{}^{}`$25) to match the instrumental point spread function. The source function, which determines the emission in each cell, can be specified as a function of the local temperature and/or density depending on the emission mechanism that is being modeled (recombination line emission, dust thermal emission, dust scattering etc.).
The physical parameters for the nebular model are well constrained by the high spatial resolution afforded by HST and NICMOS. Ground–based kinematical data that have been obtained with CSHELL at the IRTF provide radial velocity information, which imposes even stronger constraints on the nebula’s 3-D structure; these data also provide valuable information on the dynamics and evolutionary timescales of the nebula. A full presentation of these data will be made elsewhere (Kelly et al. (1999)) Note that changing the model geometry changes the shape of the image and the P-V diagram; changing the density law changes the relative location and contrast of bright and dark regions in the image and the P-V diagram, but not the shape/size of the image; finally changing the velocity law does not affect the image but changes the shape of the P-V diagram. Therefore, model images together and P-V diagrams, impose strong constraints, allowing us to develop fairly robust models for the spatial and kinematical structure of the source.
We focus here on reconstructing the 3-dimensional structure of the excited H<sub>2</sub> shell, which represents the PDR. However we also construct a model of the P$`\alpha `$ emission and compare the density and velocity structure of ionized gas that we derive, with P$`\alpha `$ and radio recombination lines respectively, and estimate a total mass for the ionized gas. Finally, we also construct a spherical, isotropically scattering dust model and compare those results with a WFPC2 continuum image. Note that each of the 3 components is modeled separately, using independent constraints provided by different images. We adopt source functions and emission coefficients appropriate for the nebular temperature and density, for each of the three components.
### 4.1 Excited H<sub>2</sub> – The Photodissociation Region
#### 4.1.1 Geometry
The H<sub>2</sub> emission is modeled as a “capped” bi-conical (hourglass-shaped) structure which shares the same axis, but lies exterior to the ionized nebula (Fig. 4). The outer dimensions of the bi-cone are well constrained by the continuum subtracted H<sub>2</sub> image (Fig. 3): The intersecting outer rings represent the maximum diameter of the cone; the major axis to minor axis ratio constrains the tilt of the cone to be $`i`$ 45°$`\pm `$5° (w.r.t. the plane of the sky). Based on an established tilt and a maximum cone diameter, the observed separation between the large rings determines the height of the cone. The radius of the waist of the hourglass is harder to constrain from the data due to projection effects. We simply assume that this dimension matches the minor axis of the ionized gas ellipsoid (§4.2), R<sub>w</sub> = 2$`{}_{}{}^{}._{}^{}`$8, though in reality, the neutral gas probably lies somewhat outside the ionized gas. The waist diameter along with the maximum diameter of the cone allows us to estimate an opening angle for the cone. There is some evidence in the images that the H<sub>2</sub> emission is not confined to a conical shell that is open at the poles, but extends all around the ionized gas ellipsoid. Wisps of H<sub>2</sub> emission are seen superposed on the northern and south-eastern parts of the ionized nebula in Fig. 2. Also, our model images show that an “open” cone cannot produce the bright rims, along the top and bottom of the structure, as are seen in the images. When a “cap” of similar thickness as the walls of the bicone is added, limb brightening along lines of sight through the caps alleviates this problem (see Figs. 5(a),(b)). By comparing intensity cuts across the model with similar cuts across the image we more accurately determine the dimensions of the H<sub>2</sub> shell. The walls of the bicone are constrained by the thickness of the observed H<sub>2</sub> rims, $`\mathrm{\Delta }`$R $``$ 0$`{}_{}{}^{}._{}^{}`$5$`\pm `$0$`{}_{}{}^{}._{}^{}`$1 (6$`\pm 1\times 10^{15}`$ cm at the adopted distance of 0.88 kpc). For example, increasing the width of the shell reduces the gap between the rims of the H<sub>2</sub> bicone and the inner, ionized gas ellipsoid; decreasing the width of the shell produces narrower looking rims than are seen in the NICMOS-only color composite image (Fig. 2(a)). Because the nebula is inclined to the line of sight, it is difficult to ascertain variations in the thickness of the H<sub>2</sub> shell as a function of latitude (particularly at low latitudes). We assume therefore that the shell has a constant width.
##### Limitations –
The cone/cap model here is an approximation to the nebula, which probably has a smoother, bi–lobed shape. However, building up the model in separate components allows us to see the effects of the cap (or lack thereof) on the P-V diagram and images. Though we do not believe that a smoother hourglass figure model would yield significantly different results, it may match the brightness in the interior (equatorial regions) of the H<sub>2</sub> images somewhat better than our current models, which appear to overestimate the H<sub>2</sub> brightness in this region.
#### 4.1.2 Kinematics
Long-slit CSHELL spectra of the H<sub>2</sub> $`v=10`$ S(1) line are invaluable in constraining the geometry of the source more tightly. The P-V diagram obtained by placing the slit along the major axis of the source is shown in Fig. 6. The emission region has an elongated elliptical shape suggesting that there are distinct red and blue–shifted components of emission both in the top and bottom halves of the nebula. The shape of the P-V diagram confirms that the nebula is oriented such that the N-E part of the shell is tilted towards the observer and the S-W part is tilted away from the observer. This orientation is consistent with that of the ionized gas ellipsoid observed by Bains et al. (1997), who used an echelle spectrograph with a velocity resolution of 6 km s<sup>-1</sup>. The four bright “spots” of emission, located symmetrically about the center of the nebula, indicate locally bright regions along the slit separated in velocity. The two bright regions closer to the center mark the position of the limb–brightened inner rims of the bicone on the slit, while the bright regions at the top and bottom correspond to the limb-brightened “caps” of the bicone, which lie almost in the plane of the sky. It is the position-velocity diagram that uniquely constrains the distribution of H<sub>2</sub> along the poles. Figs. 5 (c) and (d) show model P-V diagrams obtained with and without a cap. Clearly a conical structure with the spherical cap gives a more complete fit to the observed P-V diagram than one without the cap. A constant expansion velocity of 20 km s<sup>-1</sup> is able to reproduce the observed spectra quite well though we cannot rule out deviations from this law on smaller spatial scales. Figs. 6 (b), (c) show how the shape of the model P-V diagram changes with the velocity law. In general, a larger velocity law exponent results in a larger “kink” in the position velocity diagram, corresponding to more abrupt change in the projected velocity from the rim of the cone (large R) to the interior of the cone. A comparison with the data allows us to rule out a velocity exponent, $`\alpha 1`$. Thus, for our adopted model geometry we can rule out a constant dynamical timescale for all points in the H<sub>2</sub> shell, as one might expect for an aspherical shell expanding in a self-similar way. The velocity law implies that at a latitude of 25° the walls have a dynamical timescale of 950 yr and further out, at 36° (where the bicone ends and the spherical cap begins) the timescale is $``$ 1,300 yr. We note that since this is an expanding photodissociation front, the dynamical timescale is not the most appropriate one for characterizing the evolution (see §6.2).
#### 4.1.3 Density
Though the excitation of H<sub>2</sub> is primarily via absorption of UV photons (i.e. non-thermal, see §5.2) we assume here that the H<sub>2</sub> is thermally excited and that the level populations are described by a vibrational excitation temperature in the range $`T_{vib}6,0009,000`$ K (see Hora, Latter, & Deutsch (1999)). Later in this paper we show that this assumption gives reasonable results when compared to more detailed models of dense PDRs. If the H<sub>2</sub> emission is optically thin, then the model intensity at each pixel is simply proportional to the density of H<sub>2</sub> molecules integrated along a given line of sight:
$$ϵ_{H_2}=j_{H_2}n_{H_2}𝑑l$$
(1)
where the emission coefficient (in erg s<sup>-1</sup> cm<sup>-2</sup> per molecule) is
$$ϵ_{H_2}=\frac{A_{u,l}g_N(2J+1)h\nu }{4\pi Q(T)e^{\frac{E_u}{T}}}$$
(2)
and $`dl`$ is the spatial increment along the line of sight. The spontaneous radiative decay coefficient for the 1–0 S(1) line, $`A_{u,l}`$ = 3.47$`\times `$10<sup>-7</sup> s<sup>-1</sup> (Wolniewiez, Simbotin & Dalgarno (1998)), the nuclear spin statistical weight g<sub>N</sub> = 3 (ortho-H<sub>2</sub>) and J=3. The energy of the upper level of the transition, E<sub>u</sub> = 6947 K. Using the polynomial coefficients for the partition function, $`Q(T)`$, given in Sauval & Tatum (1984) we calculate the change in the product, $`\delta `$ $``$ Q(T) e$`^{\frac{E_u}{T}}`$ with temperature in the range 2000 K $``$ T $``$ 11,000 K. We find that $`\delta `$ changes by about a factor of 1.5 over the range 6,000 $``$ T $``$ 9,000 K. We simply assume $`\delta `$ = 220, and use this in the the above equation.
Using the equation above to fit the observed integrated intensity of H<sub>2</sub> (3.8$`\times `$10<sup>-12</sup> erg s<sup>-1</sup> cm<sup>-2</sup>; §3.1.1) and the H<sub>2</sub> radial brightness profile (Fig. 7(a)) allows us to estimate an excited H<sub>2</sub> density of $`n`$ 10–13 cm<sup>-3</sup> (depending upon the poorly known attenuation value) for a filling factor of unity. Summing over the entire nebula then yields a total mass of $`M_{ex}10^5`$ M for the excited H<sub>2</sub>. A constant density H<sub>2</sub> law provides adequate fits to the images and the P-V diagrams.
### 4.2 Ionized Gas Emission and Scattered Starlight
##### P$`\alpha `$ emission –
The ionized nebula is well reproduced by a prolate ellipsoidal shell, as has been shown previously (Atherton et al. (1979); Roelfsema et al. (1991)). Using the inclination of 45° estimated from the H<sub>2</sub> image, we calculate the lengths of the semi–major and semi–minor axes (a<sub>y</sub> and a<sub>x</sub> in Fig. 4). The thickness of the shell (though it has a broken, patchy appearance) is estimated to be about 0$`{}_{}{}^{}._{}^{}`$9 (1.2$`\times `$10<sup>16</sup> cm), from the radial cuts taken in orthogonal directions through the walls (see Table 2). Ignoring the contribution of electrons from heavier elements (i.e. assuming n<sub>e</sub> = n<sub>H</sub>), the model surface brightness in P$`\alpha `$ is given by:
$$I_{P\alpha }=j_{P\alpha }n_e^2𝑑l$$
(3)
where j (P$`\alpha `$ emission coefficient) = 3.27$`\times `$10<sup>-27</sup> erg cm<sup>3</sup> s<sup>-1</sup> ster<sup>-1</sup> assuming case B recombination, T<sub>e</sub> = 10<sup>4</sup> K (Table 4.4; Osterbrock (1989)) and l is the path length through the nebula at any given position on the projected surface of the nebula.
The density law for the ionized gas is constrained by comparing the model image (and radial cuts) with the P$`\alpha `$ image, although the clumpy nature of the ionized shell makes it somewhat difficult to fit a well–defined power law. Our results indicate that models with a constant density or slowly decreasing radial power laws (R, 0 $`\alpha `$ 0.5) models provides an adequate fit to the data (Figs. 7(b)). In models with $`\alpha `$ $``$ 1 the polar densities are considerably lower than those along the equator (Figs. 8(c),(e)); consequently radial cuts along the poles show a smaller rise in intensity going from the center outwards, than do radial cuts along the equator. The cuts through the observed image clearly do not show such a difference between the orthogonal radial cuts.
We are able to match the observed P$`\alpha `$ image with a constant density model where the density in the shell is n$`{}_{H^+}{}^{}=6\times `$10<sup>4</sup> cm<sup>-3</sup>. We also impose an upper limit to the density inside the ellipsoidal shell, n<sub>e,interior</sub> $``$ 5$`\times 10^3`$ cm<sup>-3</sup>. These densities, together with the derived geometry imply that the total ionized mass is $`M_{ion}`$ 0.018 M. (The mass might be about 50% higher if the attenuation factor is, A<sub>V</sub> = 2.97 mag – see Section 5.1). This low density inner cavity implies high temperatures ($`1\times 10^5`$ K), possibly shock heated in interacting winds. This suggests that diffuse X-ray emission from a hot gas might be detectable from NGC 7027, in addition to point source X-ray emission from the central star.
Our derived shell density agrees very well with the density derived by Roelfsema et al. (1991; 6$`\pm 0.6\times 10^4`$ cm<sup>-3</sup>) from radio continuum observations. We are unable to reproduce the total ionized mass of 0.05 M that they infer, given their ellipsoid dimensions and density (at an assumed distance of 1 kpc). Using their parameters we find an ionized mass, M$`{}_{i}{}^{}`$ 0.025 M. The PV diagrams presented in Fig. 8 (which assume that velocity is constant) are consistent with that presented for ionized gas by Roelfsema et al. (1991). A careful examination of the intensity contrast in the constant–density model suggests that this model most closely matches these data, given the S/N of the data.
##### Dust scattering –
An examination of the WFPC2 F555W image of NGC 7027 shows that the ionized nebula is surrounded by an extended region of lower surface brightness. The radial profiles indicate that the intensity falls off as I $``$ R<sup>-3</sup>, the expected index for singly scattered starlight from an optically thin region of the dust shell. We model the scattered starlight as single (isotropic) scattering from a spherical shell of 0.1 µm sized amorphous carbon grains, with a scattering cross-section, K<sub>s</sub> = 8$`\times `$10<sup>4</sup> cm<sup>2</sup> gm<sup>-1</sup> at 0.55 µm (Martin & Rogers (1987)). The photon source is assumed to be a 200,000 K black-body. Assuming an AGB dust mass loss rate of $`10^6`$ M yr<sup>-1</sup> and an expansion velocity of 10 km s<sup>-1</sup> yields a model flux density of $``$ 4.6$`\times `$10<sup>-4</sup> Jy arcsec<sup>-2</sup> at a distance of 5″ from the center of the nebula. The model radial brightness profile appears to match the observed surface brightness in the F555W image, along the minor axis of the nebula (Fig. 7c). The dust mass loss rate agrees well with the gas mass loss rate of 1.5$`\times `$10<sup>-4</sup> M derived by Jaminet et al. (1991) if the gas to dust ratio is $``$ 150.
We have tried to fit the broad–band near–IR emission with the same scattering model as used above for the WFPC2 image. We find that assuming $`\sigma _s`$ $``$ $`\lambda ^4`$ (Rayleigh scattering, appropriate for small grains) yields near–IR intensities that are considerably lower than those observed in the F160W and F205W filters. A better estimate is obtained if the scattering cross–sections goes as $`\lambda ^1`$, implying that the dust scattering is dominated by larger (“fluffy”) grains. Also, the radial profiles in the near-IR filters do not show a well–defined radial power law as seen in the WFPC2 image, suggesting that the optically thin, singly scattered starlight model is inadequate in describing the extended near-IR emission. We suspect that the emission in the near–IR broad–band filters is not simply scattered starlight but includes scattered line/f-f emission as well as possibly some dust thermal emission. A quantitative description of this “excess” near–IR emission will require a detailed model of the gas and dust emission and is not addressed further this paper.
## 5 Central Star and Nebular Excitation
### 5.1 Properties of the Central Star
The central star (CS) of NGC 7027 has been detected only with difficulty from the ground at visible wavelengths and under average seeing, because of the low contrast of the star relative to the bright nebular emission (Jacoby 1988; Heap and Hintzen 1990). NICMOS offers high resolution imaging at longer wavelengths where extinction from interstellar and circumstellar dust is much smaller than at visible wavelengths, and the contrast of the stellar point source is greatly enhanced by the high angular resolution of HST. This observational advantage allows a significant improvement in photometric data for the CS and in the stellar properties thereby inferred.
Photometry for the CS was obtained for the F110W, F160W, F190N, and F205W NICMOS filters. We have also determined CS fluxes for the F555W and F814W filters of the WFPC2 optical camera, from images in the HST data archive. Nebular photometry of the P$`\alpha `$ and H<sub>2</sub> $`v=10`$ S(1) line emission was derived from the line and continuum filter sets, F187N (line) and F190N (adjacent continuum), and F212N (line) and F215N (continuum) respectively. The photometry of the central star and the nebular line emission is summarized in Table 3.
The largest source of uncertainty in our stellar photometry arises from the spatially variable attentuation and the complex, clumpy structure of the nebula. Near the CS, this structure varies on a scale as small as the radius of the first minimum in the PSF, so background estimation and subtraction are difficult. We found that the best NICMOS photometry is achieved by subtracting a scaled P$`\alpha `$ line emission image from each of the images on which photometry is being performed. The scale factor is chosen to make the nebular background have a zero mean in the vicinity of the CS. (The P$`\alpha `$ image is obtained by subtracting the F190N continuum image from the F187N P$`\alpha `$ narrow band image, so that the central star is removed, leaving only the nebular line emission.) In this procedure we assume that the P$`\alpha `$ line emission is a good tracer of total nebular line and continuum emission in the other NICMOS filters, at least in the region near the CS. We find that this method successfully removes the spatially variable nebular emission near the CS.
In the WFPC2 images the nebular surface brightness shows smaller variations relative to the CS. In this case we employ more conventional photometry and assume that the surface brightness in an annulus immediately surrounding the CS is a fair representation of the background. The photometry was converted to flux density by applying the photometric calibration factor (PHOTFLAM) provided by STScI in the image header.
The CS photometry is dereddened using an attenuation value at the position of the CS of $`A_V=3.17`$ mag. The P$`\alpha `$ line emission photometry is corrected for an average nebular attenuation of $`A_V=2.97`$ mag (Robberto et al. 1993). The correction at each wavelength is derived from the mean extinction curve given in Table 3.1 of Whittet (1992).
We calculate a hydrogen Zanstra temperature (cf. Osterbrock 1989) for each CS photometric value in Table 1, with the assumption of a blackbody CS spectrum. For the nebular line emission, we use the integrated P$`\alpha `$ line flux given in Table 3. The greatest source of error in this calculation results from uncertainties in the nebular conditions on which the recombination coefficients depend. We adopt a nebular electron temperature of 15000 K and a nebular electron density of 10<sup>4</sup> cm<sup>-3</sup> (cf. Roelfsma et al. 1991; Masson 1989), and use the corresponding Case B total recombination coefficient, $`\alpha _B`$, and P$`\alpha `$ emissivity calculated by Storey & Hummer (1995). The derived Zanstra temperatures of the CS are compiled in Table 3.
The derived CS temperatures are remarkably consistent (see Table 3). We find $`T_{}=198,100\pm 10,500`$ K. Adopting a distance of 880$`\pm `$150 pc (Masson 1989), we find from the photometric data of Table 3 a mean CS luminosity of 7710 $`\pm `$ 860 L, and a photospheric radius of $`5.21\times 10^9`$ cm. If we place the CS on the evolutionary tracks for hydrogen-burning post-AGB stars of Blöcker (1995), the CS lies very close to the theoretical track for an initial mass of 4 M and a final core mass of 0.696 M. If this model is appropriate, the CS would have left the AGB only about 700 years ago. This timescale is consistent with the estimated age of 600 years for the ionized nebula (Masson 1989) based on the expansion velocity and size. The stellar evolutionary timescale is also comparable to the dynamical timescale of the H<sub>2</sub> lobes estimated in §4.1.2 above. Jaminet et al. (1991) have observed the neutral molecular envelope which surrounds the ionized gas and determined a mass of 1.4 (+0.8, –0.4) M. If the CS initial mass was in fact 4 M, and the final core mass is 0.7 M, then the star must have shed 3.3 M over the course of its evolution. Of this amount, 40 – 50% was lost in the last 10<sup>4</sup> years, when the mass loss rate was some $`1.5\times 10^4`$ M yr<sup>-1</sup> as inferred from models for the CO emission of the envelope (Jaminet et al. 1991). The balance of the mass loss might have occurred more slowly over a much longer time period, but the lack of clear observable signatures of this matter make it impossible to infer the rate or timescale of this earlier phase with any precision. The stellar properties found for NGC 7027 are summarized in Table 4.
### 5.2 Molecular Hydrogen Excitation and the Morphology of NGC 7027
The near-infrared spectrum of molecular hydrogen is the result of slow electric quadrapole transitions within the ground electronic state vibration-rotation levels. A possible excitation mechanism likely in PNe is via the absorption of far-ultraviolet (FUV) photons in the Lyman and Werner bands of H<sub>2</sub>. Upon absorption of a FUV photon, the molecule is left in an excited electronic state from which $``$10% of the decays out of that state result in transitions to the ground electronic state vibrational continuum and molecular dissociation occurs. Most often the molecule is left in an excited vibration-rotation level in the ground electronic state. The result is an identifiable near-IR fluorescence spectrum. This process is the primary way H<sub>2</sub> is destroyed, by photons and the peak in ro-vibrationally excited molecules defines the H<sub>2</sub> photo-dissociation region (PDR). In addition, the excitation of the near-IR H<sub>2</sub> spectrum can occur in regions with kinetic temperatures of $`T_K1000`$ K and total densities $`n_{tot}10^4`$ cm<sup>-3</sup>. In PNe such temperatures are typically thought to be associated with shockwaves of moderate magnitude and over a limited range ($`V_s550`$ km $`\mathrm{s}^1`$; with the higher end of the range for $`C`$-type shocks only; Hollenbach & Shull (1977); Draine & Roberge (1982)).
Graham et al. (1993a) argued, primarily from morphology, that the near-IR H<sub>2</sub> emission from NGC 7027 is caused by FUV excitation in the PDR. An analysis of near-IR spectra of NGC 7027 (Hora, Latter, & Deutsch (1999)) confirms that hypothesis. There is no evidence of shock excitation, though localized shocked regions cannot be ruled out.
In a PDR the H<sub>2</sub> emission is observed in a thin transition region where the H<sub>2</sub> column density of FUV excited molecules peaks. Interior to the transition region the molecules have mostly been dissociated. Because the excitation process is by absorption of line photons, the depth of excitation is strongly limited by self-shielding of the molecules to Lyman and Werner band line radiation. The total extent of the transition region depends on molecular and total density, gas to dust ratio, and the strength and shape of the radiation field.
Comparison with quasi-steady state PDR models (Latter & Tielens, in preparation) finds that an observed average H<sub>2</sub> $`v=10`$ S(1) line brightness of $`8.0\times 10^4`$ erg s<sup>-1</sup> cm<sup>-2</sup> ster<sup>-1</sup> (§3.1) is consistent with a C-rich PDR exposed to a single star of temperature 200,000 K and a moderate total density of $`n_{tot}5\times 10^4`$ cm<sup>-3</sup>. The predicted $`v=21`$ S(1) to $`v=10`$ S(1) line ratio ranges from 0.12 to 0.06 and is nearly identical to that found by analysis of the H<sub>2</sub> spectrum (Hora, Latter, & Deutsch (1999)). For a core mass of 0.696 M and the observed H<sub>2</sub> line brightness, time dependent models of PDR evolution in PN suggest a PDR lifetime of $`t1000`$ years (Natta & Hollenbach (1998)) consistent with estimates of the time since the object left the AGB (§5.1). The models are also consistent with a thin transition (H<sub>2</sub> emitting) region like that found from the 3-D models discussed in §4. The nonuniformity and patchiness of the transition region must be a result of structure in the circumstellar material itself. Instabilities of the type required are not understood in an interacting wind, photon-dominated environment. The “waves” present on the illuminated interior of the H<sub>2</sub> emission region resemble those seen in the HST/WFPC2 images of MyCn 18 (the “Etched Hourglass Nebula;” Sahai et al. (1999); Dayal et al. (1999)), and although are seen in very different components of the gas (ionized versus molecular) most likely are a result of a similar instability.
It is important to consider whether the overall shape of the PDR (as it extends from the ionized core to the biconical, or bubble shaped H<sub>2</sub> emission) as found from our HST/NICMOS images is consistent with what is expected for an evolving PDR in a distribution of material likely to be present in a post-AGB circumstellar envelope. First, we must determine what is a likely distribution of circumstellar material. It has been known for some time that post-AGB objects and PNe generally do not have a spherical (radially symmetric) distribution of gas and dust. The most extreme examples include AFGL 2688 (the “Egg Nebula”), AFGL 618, M 2-9, and M 1-16 (to list just a few; see, e.g., Sahai et al. (1998); Latter et al. (1995); Latter et al. (1992); Hora & Latter (1994) and references therein). These objects have a predominantly axially symmetric distribution of circumstellar material such that the density is greatest at the equator and lowest near the two poles. With support from general morphological studies, hydrodynamic modeling, and polarimetric observations (e.g. Balick (1994); Trammell, Dinerstein & Goodrich (1994); Balick & Frank (1997)) this type of distribution (to a widely varying degree) appears to be generally true for most, or all post-AGB objects. NGC 7027 is not an “extreme” bipolar nebula (of the type listed above), but we can assume that the equatorial to polar density contrast as the star evolved off the AGB has a functional form that goes roughly like $`\mathrm{sin}(\mathrm{\Theta })`$ with additional factors of order unity. Functions like this have been used to model the scattered light and molecular distribution from a number of bipolar PPNe and PNe (Yusef-Zadeh, Morris, & White (1984); Latter et al. (1992); Latter et al. 1993b ).
In an environment with a single, central FUV source and with a constant gas to dust ratio, the dominant parameters that will determine the radial structure of the PDR are the density structure, $`\rho (r)`$ and the distance to the central source. For guidance some new computations were made using a code developed for a related purpose, and this topic will be explored quantitatively elsewhere (Latter & Tielens, in preparation). To first order, a PDR in a lower density region will have the ionized/neutral interface farther from the central source when compared to a higher density region. Similarly, the region of strongest H<sub>2</sub> FUV excitation will be farther out and have a somewhat larger linear extent than in a higher density region. This H/H<sub>2</sub> interface region will approximately follow iso-column density contours at $`4\times 10^{21}`$ cm<sup>-2</sup> or $`A_V2`$ magnitudes. Details will depend strongly on the exact distribution of material. But, it is evident that an axially symmetric distribution of circumstellar material of the type discussed above will produce a PDR morphology like that observed in NGC 7027. More specifically, it is expected to be an ellipsoidal or bipolar H<sup>+</sup>/H interface, surrounded by a bipolar bubble structure defining the H/H<sub>2</sub> interface and revealed in FUV excited H<sub>2</sub> emission – like that modeled here. Somewhat farther out, the C/CO interface region is, as shown by the data of Graham et al. 1993a , consistent with that found for the inner regions. The morphological evolution of the H<sup>+</sup>/H interface region and the ionized core is not predicted by strict PDR evolution alone. The internal density structure appears to have equilibrated in the shell and cavity (§4.2) because the sound crossing time is much shorter than the expansion timescale. The outward motion of this interface will be faster in the lower density polar regions, magnifying the asymmetric shape. If the star is on the cooling track, then the ionization front might better be described as a “recombination” front as the number of ionizing photons goes down (Schoenberner 1999, personal communication).
Nested hourglass shapes and ellipsoids are predicted by numerical hydrodynamical simulations (such as Mellema & Frank (1995); Icke (1988)), but for interaction speeds much higher than so far seen in NGC 7027. In addition, the molecular hydrogen emission from NGC 7027 is from a strong PDR, not a shocked wind interface. There is similarity in the overall morphology of the PDR, including the ionized core in NGC 7027 to emission morphologies predicted in strong interacting winds starting in an axisymmetric density distribution. Because of this apparent similarity, a full description of the observed structure and its evolution requires coupled time-dependent PDR and hydrodynamic modeling. When combined with detailed chemical modeling of the neutral regions, the data presented here can be used to determine the exact density distribution of the circumstellar envelope.
## 6 Discussion
### 6.1 Previous Kinematic Results
In a detailed multiwavelength study, Graham et al. (1993a) found that uniform, radial expansion of a prolate spheroid, tilted to the line of sight is consistent with the millimeter-wave (CO) position-velocity maps. However the CO P-V diagram differs from the H<sub>2</sub> P-V diagram presented here. This is not unexpected since the observed emission from the two molecules is produced under very different excitation conditions. Therefore the two molecules trace different spatial/kinematical structures. Observations have been made with the “BEAR” instrument (a 256$`\times `$256 pixel HgCdTe camera coupled to a Fourier transform spectrometer, with a spectral resolution $``$ 52.2 km $`\mathrm{s}^1`$) on the Canada-France-Hawaii Telescope (Cox et al. (1997)). Cox et al. found that the velocity field appeared to them as an equatorial “torus” suggests that north pole is tilted away from observer. This result is inconsistent with our interpretation of the HST images, our 3-D modeling of those images, and the higher spectral resolution CSHELL data. The high spatial resolution and sensitivity of the HST images shows that the front part of the top rim lies in front of the the ionized nebula (not behind), suggesting that the top (northern side) is tilted $`toward`$ the observer. We cannot resolve this discrepancy with the data available. We suggest that the spectral resolution available to Cox et al. (1997) was not sufficent to fully deconvolve the subtle, but complex velocity structure that must be present (Kelly et al. (1999)). In addition, the model schematic given by Jaminet et al. (1991) shows the same orientation as Cox et al. (north pole away) for the fast molecular wind traced by CO $`J=32`$ line emission, which presumably has a similar flow direction as the ionized gas. It is entirely possible that the flow along the back of the cone (or bubble) has been confused with a flow along the polar axis, which in our model would not be the dominant emitting region. Bains et al. (1997) found from their high spectral resolution radio and optical observations, an orientation for NGC 7027 that is consistent with ours: i.e. that the NE lobe is blue-shifted, SW lobe is red-shifted.
### 6.2 Evolution of the PDR in NGC 7027
A full modeling of PDR evolution is beyond the scope of this paper. We can, however, make an assessment of such PDR evolution based on available models. The morphological evolution of the PDR will depend most strongly on the distribution of material in the circumstellar envelope. For NGC 7027 the post-AGB circumstellar envelope must have been axially symmetric with a gradual decrease in density from the equatorial region to the poles. A dense equatorial disk cannot be ruled out, but there is no evidence in these data, or other molecular data (Bieging, Wilner, & Thronson (1991); Jaminet et al. (1991); Graham et al. 1993a ), that requires or even suggests one. For that reason, we consider such a structure as unlikely to be present. The density distribution appears typical for other post-AGB objects with an equatorial to polar density contrast of $`210`$, such that it is not as great as the extreme bipolar proto-PNe (see §5.2). In much less than the cooling timescale of the central star, the PDR will be established quickly in all directions, with the fastest evolution in the polar regions. The result will most likely be for NGC 7027 to become a butterfly-type nebula similar in gross properties to (e.g.) NGC 2346. Wind interactions for this type of morphological evolution are not required.
The timescale for PDR evolution is highly uncertain (other timescales are discussed in §5.1), but reasonable estimates can, and have been made (see, e.g., Tielens (1993)). It is straightforward to show that for a constant mass loss rate and parameters typical of high mass loss rate carbon stars, the timescale for a PDR to completely move through the circumstellar envelope can be estimated by integrating the equations of molecular formation and destruction with respect to radius and time, as was done by Tielens (1993). For a constant wind velocity $`V_w=20`$ km $`\mathrm{s}^1`$, a constant mass loss rate $`\dot{M}=10^4`$ M yr<sup>-1</sup>, and a stellar luminosity $`L_{}=6000`$ L, it can be shown that the dissociation front travels a distance $`r_i3\times 10^{16}`$ cm in $`t_i65\mathrm{years}`$ with a PDR speed of $`V_i263\mathrm{km}\mathrm{s}^1`$ (see Tielens (1993) for a more detailed discussion; see also Natta & Hollenbach (1998)). Tielens suggested clumping as how the evolution is slowed in real nebulae. Although widespread clumping and structure is seen in these data, it does not appear at the size scales required, or marginally so at best ($`R10^{16}`$ cm and $`A_V4`$ mag). Even if clumping is important in localized regions, it is more likely that we now know several of the assumptions used to derive this evolutionary timescale are not valid, and if properly accounted for will tend to increase the timescale. The mass loss rate is not constant. The “ring” structure clearly visible in AFGL 2688 (Latter et al. 1993b ; Sahai et al. (1998)) indicates a time varying mass loss rate on a timescale of hundreds of years. The same type of structures are seen in other objects observed (Kwok et al. (1998)), including NGC 7027 (see, e.g., Bond et al. (1997)). In addition, the distribution of material is not spherical. The degree to which the mass loss rate is varying is not known. Nor is how a non-spherical distribution of attenuating and shielding material will alter the evolutionary timescales (see, however, Natta & Hollenbach (1998)). Additional work must be done to characterize these properties and how they relate to evolution. From the above expression and an assessment of the impact of uncertain parameters, it seems likely that the timescale for NGC 7027 to evolve away from its current state will be only a few hundred to a few thousand years (see also Natta & Hollenbach (1998)). Wind interactions will complicate the picture by adding hydrodynamic effects and changes to the density structure. However, evolution to the current epoch appears to be described fully by the interaction of FUV and ionizing photons with an AGB wind that has (or had) a somewhat enhanced mass loss rate in the equatorial plane. While there is evidence for recent jet interactions with the outflow, the jets have not dominated the morphological shaping. It is worth monitoring NGC 7027 at very high spatial resolution for morphological changes caused by jets and UV photons over the next several decades.
## 7 Summary
The HST/NICMOS observations presented in this paper reveal the detailed structure and morphology of the young planetary nebula NGC 7027 with unprecedented clarity. The molecular photodissociation region is found to be of an apparently different overall structure from that of the ionized core. We have constructed 3-dimensional, axisymmetric models for the morphology of NGC 7027 based on these NICMOS near-infrared data and HST archive visible light data. The object can be well described by three distinct components: an ellipsoidal shell that is the ionized core, a bipolar hourglass structure outside the ionized core that is the excited molecular hydrogen region (or the photodissociation region), and a nearly spherical outer region seen in dust scattered light and is the cool, neutral molecular envelope. It is argued that such a structure is consistent with one which would be produced by a photodissociation region with a central source of far-ultraviolet photons surrounded by an axisymmetric distribution of gas and dust that has a decreasing total density with increasing angle toward the polar regions of the system.
The central star is clearly revealed by these data. We are able to determine the most accurate value for the temperature of the star $`T_{}=198,100\pm 10,500`$ K with a luminosity of $`L_{}7,700`$ L. From various indicators, including the central star timescale, dynamical timescale, and timescale for evolution of the photodissocation region, we conclude that NGC 7027 left the asymptotic giant branch 700 to 1000 years ago. For such rapid evolution, it is likely that objects like NGC 7027 do not go through a proto-planetary nebula phase that is longer than a few tens of years, and as such would be nearly undetectable during such a brief phase.
There is strong evidence for one, and possibly two highly collimated bipolar jets in NGC 7027. The jets themselves are not seen with the present data, but the disturbance caused by them is clearly visible. It is possible that this jet, or jets have shut off during the current epoch of evolution. Such jets have been seen in other planetary nebulae, but NGC 7027 might be the youngest in terms of evolution. There are nonuniformities and wave-like structures seen in these data that must be caused by poorly understood instabilities in the outflowing wind.
We conclude that while wind interactions will be important to the future evolution of NGC 7027, it is the evolution of the PDR, as the central star dissociates, then ionizes, the circumstellar medium, that will dominate the apparent morphological evolution of this object. This is in contrast to the idea that interacting winds always dominate the evolution and shaping of planetary nebulae. Evolution driven by far-ultraviolet photons might be more common in planetary nebulae than has been previously thought (see, e.g., Hora, Latter, & Deutsch (1999)). Such photon-driven evolution will be more important in objects with higher mass central stars (and therefore high luminosity, temperature, and UV flux), than in those with low mass stars. But, it is the high mass objects that also have the higher density molecular envelopes ejected during a high mass loss rate phase on the asymptotic giant branch, and as such will show the strongest emission from a photodissociation region. NGC 7027 is perhaps the most extreme example known at this time, but it is not unique – except for the important and very brief moment in evolution at which we have found it.
This work is dedicated to the memory of Christopher Skinner, our Contact Scientist at STScI, our colleague, and our friend. We miss his support, comments, and criticism. He challenged us to search for the truth, and through that he enriched us all. We gratefully thank the members of the NICMOS Instrument Development Team at the University of Arizona for help, suggestions, comments, and data; especially Marcia Rieke, Dean Hines, and Rodger Thompson. We also thank Sun Kwok and Nancy Silbermann for useful conversations, and Tom Soifer for reading the manuscript. Support for this work was provided by NASA through grant number GO-7365-01 from the Space Telescope Science Institute, which is operated by AURA, Inc., under NASA contract NAS5-26555. W.B.L. and J.L.H. acknowledge additional support from NASA grant 399-20-61 from the Long Term Space Astrophysics Program. |
warning/0002/hep-lat0002013.html | ar5iv | text | # 1 Introduction
## 1 Introduction
The goal of lattice QCD is the computation of physical quantities such as hadron masses and matrix elements using numerical Monte Carlo methods. This has proved to be an ambitious programme because after discretisation of the path integral and generation of a sufficiently large number of independent configurations several limits must be considered:
1. The box size. This is currently at $`1.53\text{fm}`$, and should be compared with the nucleon rms radius of $`0.8\text{fm}`$.
2. The chiral limit $`m_q0`$. The $`u/d`$ and $`s`$ quarks are light quarks.
3. Continuum limit $`a^k0`$ ($`k=2`$ if we choose $`O(a)`$ Symanzik improved fermions, staggered fermions or Ginsparg-Wilson fermions; for Wilson fermions we expect the discretisation effects to have $`k=1`$).
If all these limits can be successfully taken then presumably QCD will reproduce nature. Although first attempts in this direction are being made, , it will require much faster computers to achieve this goal. To reduce the computational effort often the quenched approximation is employed when the fermion determinant is simply set to a constant. However then new problems arise (or are exacerbated):
1. Spurious quenched chiral logarithms appear as $`m_q0`$.
2. The appearance of exceptional configurations.
3. Consistency of the final results when comparing with their experimental or phenomenological values.
In this talk we shall consider points 1 and 2 numerically using $`O(a)`$ improved fermions.
## 2 Chiral perturbation theory and quenched chiral logarithms
This was developed by Bernard and Golterman, and Sharpe, . What do we expect? The quenched pseudoscalar effective chiral Lagrangian gives
$`(am_{ps})^2`$ $``$ $`(a\stackrel{~}{m}_q)^{\frac{1}{1+\delta }},`$
$`a^2g_P`$ $``$ $`0|\widehat{𝒫}|ps[(am_{ps})^2]^\delta ,`$
$`am_{ps}af_{ps}`$ $``$ $`0|\widehat{𝒜}_4|psam_{ps}.`$ (1)
As the $`\eta ^{}`$ remains light and has a single and double pole in its propagator, this latter term at $`p^2=0`$, $`(m_0^2/m_{ps}^2)(1/m_{ps}^2)`$ acts like an extra vertex giving a singular correction to the usually harmless loop term $`m_{ps}^2\mathrm{ln}(m_{ps}/\mathrm{\Lambda }^2)`$ of $`m_0^2\mathrm{ln}(m_{ps}/\mathrm{\Lambda })^2`$, so that the logarithmic term becomes singular, . This can then be summed to give eq. (1). Normally PCAC, $`𝒜=2\stackrel{~}{m}_q𝒫`$, would give us $`(am_{ps})^2a\stackrel{~}{m}_q`$. Thus the non-zero $`\delta `$ leads to singular behaviour for $`m_{ps}`$ and $`g_P`$ in the chiral limit for quenched QCD. Simple estimates lead to an expectation for $`\delta 0.10.2`$. $`\mathrm{\Lambda }`$ is a cut-off on the $`\eta ^{}`$ loop so $`\mathrm{\Lambda }\text{ }<900\text{MeV}`$ and above this scale any singularities are surely damped out. For consistency, we would expect little difference between the PCAC quark mass, $`a\stackrel{~}{m}_q`$, and the ‘standard’ quark mass, $`am_q\frac{1}{2}(1/\kappa 1/\kappa _c)`$.
Similarly by considering the vector/baryon effective chiral Lagrangians, it can be shown that
$$am_{V,N}=c_0+c_1am_{ps}+c_2(am_{ps})^2+O((am_{ps})^3),$$
(2)
where there is an extra linear term present, as the $`\eta ^{}`$ gives to the usual term $`(am_{ps})^3`$ an additional term $`m_0^2am_{ps}`$.
## 3 Numerical results
How do these theoretical considerations fare with the numerical data? Numerically, it advantageous to use the PCAC quark mass, $`\stackrel{~}{m}_q`$ because this quark mass can be found very accurately and does not depend on the first-to-be-determined parameter $`\kappa _c`$. It is convenient to consider the ratio $`a\stackrel{~}{m}_q/(am_{ps})^2`$. This is shown in Fig. 1, , for degenerate quark masses
at $`\beta `$ values of $`6.0`$, $`6.2`$ and $`6.4`$. To give an idea of scales, we note that using the $`r_0`$ ‘force scale’ then $`m_\pi `$ lies almost at the chiral limit (within our numerically accuracy there is no difference between these points) and a hypothetical pseudoscalar meson, composed of the strange quark and its antiquark, lies at about $`(r_0m_{ps})^23.13`$ (when using $`m_K`$). For the charm quark (using $`m_D`$) we find $`(r_0m_{ps})^244.9`$, way off the plot scale.
It is to be seen from the picture that from the strange quark mass to heavier quark masses we have linear behaviour. (Indeed the linearity seems to hold until rather heavy quark masses, say $`m_q\text{ }<\frac{1}{3}m_c`$.) Below the strange quark mass, there seems to be a (sharp?) break in this behaviour. Possibly we can attribute this to the onset of quenched chiral logarithms. $`m_{ps}=\sqrt{2}m_K`$ corresponds here to about $`700\text{MeV}\mathrm{\Lambda }`$ so one might hope that any quenching effects are suppressed above this value. We now check the behaviour between $`a\stackrel{~}{m}_q`$ and $`am_q`$. In Fig. 2 we plot $`a\stackrel{~}{m}_q/am_q`$ against $`am_q`$.
While there seems to be a reasonably linear relation between the two (lattice) definitions of the quark mass above $`am_s`$, below we again see deviations. This fit is somewhat sensitive to the value of $`\kappa _c`$ used; although the general picture shown in Fig. 2 never seems to change significantly. (Indeed using all the light quark data in the fit still produces a similar result.) This perhaps obscures the interpretation of Fig. 1 as being due to quenched chiral logarithms. Nevetheless, due to problems in determining $`am_q`$, we prefer to use the results with $`a\stackrel{~}{m}_q`$, .
To try to expose the small quark region, and to determine $`\delta `$ (if the deviations are due to quenched chiral logarithms) then it is convenient to plot the logarithm of Fig. 1. This is done in Fig. 3,
where we expect the slope to be $`\delta `$ (for the fitted quarks masses below the strange quark mass). We find for $`\beta =6.0`$, $`\delta 0.12(4)`$, and for $`\beta =6.2`$, $`\delta 0.06(2)`$ (for $`\beta =6.4`$ there is not enough data). For $`\beta =6.0`$, at least, there is reasonable agreement with the theoretical prejudice; for $`\beta =6.2`$ the value seems small.
Despite the above results it should be noted, as discussed above, that it is notoriously difficult to numerically detect quenched chiral logarithms. Indeed the above effects may simply be due to finite-size effects or ‘exceptional configuration’ problems. We have only been able to check very few quark mass points for finite size effects. The impression is that they are small; this is backed up by , who work on a larger lattice.
Exceptional configurations are seen as either the non-convergence of the fermion matrix inversion or the correlator seems to have a ‘fake source’ at some $`t`$ value. The problem is more severe for $`O(a)`$ improved fermions than for Wilson fermions and increases as $`\beta `$ and/or $`c_{sw}`$ and/or $`m_q`$. (This is the main obstacle for the $`O(a)`$ improvement programme going below about $`\beta 6.0`$.) The reason for this problem is due to the presence of small real eigenvalues in the fermion matrix. In an experiment, , we have chirally rotated the lattice quark mass, $`am_q`$, away from the real axis; the same configuration then gave a well behaved pion propagator. (We might then expect some mixing in the correlation function of the particle with its parity partner, however for the pion in the quark model this partner does not exist.) So perhaps simply throwing away the configuration does not affect the spectrum (?). It is desirable to have a crude indicator of whether we have an exceptional configuration. In the simple proposal was made to look at the pion norm,
$$\mathrm{\Pi }(\{U\})=\underset{\stackrel{}{x},t}{}|\gamma _5G(\stackrel{}{x},t;\stackrel{}{0},0;\{U\})\gamma _5|^2.$$
(3)
(In our application, the source $`(\stackrel{}{0},0)`$ was also Jacobi smeared.) In Fig. 4,
we show a sequence of pion norms for $`\beta =6.2`$. To decide on a criterion for an exceptional configuration (ie a spike in the pion norm) is not so easy. Some are obvious, for example from the pictures we have at $`\kappa =0.1354`$, $`n_{conf}=217`$ a problem. (The corresponding pion propagator is shown in .) Closer to the critical point than here more spikes have been seen in Wilson data, . To be safe, we have actually chosen a more conservative local criterion where if at any $`t`$ value the (pion) correlation function fluctuates more than $`5`$ standard deviations from the local average we reject the configuration. This leads at $`\beta =6.2`$, $`32^3\times 64`$ for $`\kappa =0.1352`$, $`0.1354`$, $`0.13555`$ to rejection rates of about $`2`$, $`4`$ and $`6\%`$ respectively. (For the latter two $`\kappa `$ values about $`15\%`$ and $`33\%`$ of these rates were actually due to non-convergence of the inverter.) So in conclusion: while we feel that for any lighter quark mass than those considered here exceptional configurations become a real problem, at our masses while they are a nuisance, they do not distort the numerical result.
We now turn to a consideration of the decay constant $`f_{ps}`$. From eq. (1) we expect that $`af_{ps}`$ has no quenched chiral singularity in it, while $`a^2g_P`$ diverges in the chiral limit. In Fig. 5
we plot the unrenormalised $`af_{ps}`$. The results seem smooth over the whole quark mass range, with no singular behaviour. Looking at the ratio $`af_{ps}/a^2g_P`$ we see the same behaviour as in Fig. 1, ie for smaller quark masses than the strange quark mass a bend is seen in the data. As $`af_{ps}/a^2g_P[(am_{ps})^2]^\delta `$ then taking the logarithm gives a direct estimate of $`\delta `$. In Fig. 6
we show this, with fit values $`\delta 0.10(3)`$ ($`\beta =6.0`$) and $`\delta 0.05(2)`$ ($`\beta =6.2`$) consistent with the previous results.
Finally we consider chiral extrapolations of the nucleon and rho masses. In Figs. 7, 8
we show the results, together with a phenomenological fit
$$(am_{\rho ,N})^2=b_0+b_2(am_{ps})^2+b_3(am_{ps})^3.$$
(4)
In distinction to eq. (2) this does not have a quenched linear chiral term. (As there is curvature in the results above the strange quark mass, we have considered $`(am_{\rho ,N})^2`$ rather than $`am_{\rho ,N}`$ and included a cubic term in eq. (4). This gave a better fit function for the data.) While, for the nucleon this gives a good description of the data over the whole quark mass range and it is thus difficult to say in this case whether a linear term is necessary or not, the $`\rho `$ data might be showing some deviations for small quark masses.
## 4 Conclusions
Our main conclusion is that in quenched QCD there seems to be a dangerous region for quark masses $`m_q\text{ }<m_s`$. If we are interested in the strange quark mass or particles such as $`m_K`$, $`m_K^{}`$, $`\mathrm{}`$, or decay constants such as $`f_K`$, $`f_K^{}`$, $`\mathrm{}`$, this does not represent a problem. For quantities involving only the $`u`$ and $`d`$ quarks, it is probably best to adopt the pragmatic approach of making fits for $`m_q\text{ }>m_s`$ and then to extrapolate this to $`m_qm_{u/d}`$, ie to chiral limit. (See, for example, the results for the quark mass, .) Evidence for chiral quenched logarithms is mixed – the best signal seems to be for the pion and its associated decay constant. Other channels seem to be less unambiguous. Indeed, as detecting and measuring quenching effects can be quite difficult this would indicate that the quenched approximation often seems to be working quite well.
The above results should be regarded as preliminary. We hope to present full results shortly, , including continuum extrapolations (considered in the talk, but not described here).
## Acknowledgements
The numerical calculations were performed on the Quadrics QH2 at DESY (Zeuthen) as well as the Cray T3E at ZIB (Berlin) and the Cray T3E at NIC (Jülich). We wish to thank all institutions for their support. |
warning/0002/hep-ph0002272.html | ar5iv | text | # A New Source for Electroweak Baryogenesis in the MSSM
\[
## Abstract
One of the most experimentally testable explanations for the origin of the baryon asymmetry of the universe is that it was created during the electroweak phase transition, in the minimal supersymmetric standard model. Previous efforts have focused on the current for the difference of the two Higgsino fields, $`H_1H_2`$, as the source of biasing sphalerons to create the baryon asymmetry. We point out that the current for the orthogonal linear combination, $`H_1+H_2`$, is larger by several orders of magnitude. Although this increases the efficiency of electroweak baryogenesis, we nevertheless find that large CP-violating angles $`0.15`$ are required to get a large enough baryon asymmetry.
PACS: 98.80.Cq
\]
1. A highly constrained proposal for the origin of baryonic matter in the universe is electroweak baryogenesis in the minimal supersymmetric standard model (MSSM) . Unlike many other baryogenesis mechanisms, this one has strong prospects for being falsified in upcoming experiments, due to its need for some light exotic particles, notably a right-handed top squark which is lighter than the top quark .
The basic mechanism is intuitively clear: particles interact in a CP-violating manner with bubble walls, which form during the first order electroweak phase transition, when the temperature of the universe was near $`T=100`$ GeV. This causes a buildup of a left-handed quark density in excess of that of the corresponding antiquarks, and an equal and opposite right-handed asymmetry, so that there is initially no net baryon number. The left-handed quark asymmetry biases anomalous sphaleron interactions, present within the standard model, to violate baryon number preferentially to create a net quark density. The resulting baryon asymmetry of the universe (BAU) soon falls inside the interiors of the expanding bubbles, where the sphaleron interactions are shut off, and thus baryon number is safe from subsequent sphaleron-induced relaxation to zero.
Despite the simplicity of this idea, a quantitatively accurate calculation is difficult to carry out. Various approximations have been made, leading to a variety of formalisms which give somewhat conflicting predictions for the dependence of the BAU on the parameters of the MSSM. Although most authors agree on the diffusion equations which govern the generation of the left-handed quark asymmetry, there is less consensus about how to derive the source terms which appear in these equations.
In a previous paper , we advocated an approach based on the classical force on particles due to their spatially varying masses as inside the bubble wall. It is straightforward to solve for this force, put it into the Boltzmann equations, and derive corresponding diffusion equations. No ad hoc assumptions are needed; one only requires that the width of the wall be significantly larger than thermal de Broglie wavelengths of particles in the plasma, to justify an expansion in derivatives of the background Higgs field, which constitutes the bubble wall. Detailed calculations of the wall width confirm that this is indeed a good approximation .
One puzzling conflict between our earlier work and that of others was that we derived a source which remains large when the ratio of the two Higgs fields, $`H_2(x)/H_1(x)`$, is constant inside the bubble wall. Other authors found sources proportional to the derivative of this quantity. Careful analyses of the shape of the wall have shown that in fact $`H_2/H_1`$ remains constant to within a part in $`10^210^3`$ , so that the dependence on $`d(H_2/H_1)/dx`$, if correct, would result in a large suppression of the generated BAU.
In this Letter we explain the origin of the apparent discrepancy: it is the result of different choices about which linear combination of Higgsino currents is used to source the diffusion equations. References considered only the combination $`H_1H_2`$, based on the observation that the orthogonal source, $`H_1+H_2`$, is driven to zero in the limit that interaction rate $`\mathrm{\Gamma }`$ proportional to the top quark Yukawa couplings becomes infinite. Although this approximation is convenient because it simplifies the network of coupled diffusion equations, in the present application it can lead to a serious underestimate of the BAU. Parametrically, the $`H_1+H_2`$ source is suppressed by $`(\mathrm{\Gamma }D_{\stackrel{~}{h}})^{1/2}`$, where $`D_{\stackrel{~}{h}}20/T`$ is the Higgsino diffusion constant . However for realistic values, $`(\mathrm{\Gamma }D_h)^{1/2}1`$, so there is no actual suppression, eliminating the need for large CP violating phases to get the observed BAU, which we had found with the $`H_1H_2`$ source . Here we will demonstrate how this comes about, while at the same time updating our earlier computation so as to give a quantitatively accurate determination of the BAU as a function of relevant parameters in the MSSM.
2. To explain the how the classical force mechanism works, we first review the simple case of a top quark with a spatially varying complex mass , $`m(x)=yH(x)e^{i\theta (x)}`$. By solving the Dirac equation $`(i/|m|\mathrm{cos}\theta i|m|\mathrm{sin}\theta \gamma _5)\psi `$ to first order in derivatives (the WKB approximation), one finds that a particle of energy $`E`$ experiences a spin ($`s`$) dependent force
$$F=\frac{dp}{dt}=\frac{(m^2)^{}}{2E}+\frac{s}{2E^2}\left(m^2\theta ^{}\right)^{}.$$
(1)
The spin dependent part of the force has opposite sign for antiparticles, which causes the distributions of like-helicity fermions and antifermions to be separated in the vicinity of the wall. For relativistic particles, we can approximate helicity by chirality, and speak of spatially varying chemical potentials, $`\mu _{L,R}`$ for the left- and right-handed components; these are related to the asymmetry between the particle and antiparticle densities by $`n(x)\overline{n}(x)=\mu (x)T^2/6`$. Diffusion equations can be derived by inserting the force (1) into the Boltzmann equation, doing an expansion in moments of the distribution functions, and truncating the expansion. They have the form
$$D\mu _{L}^{}{}_{i}{}^{\prime \prime }v_w\mu _{L}^{}{}_{i}{}^{}+\mathrm{\Gamma }\mu _{L}^{}{}_{i}{}^{}=S_i(x),$$
(2)
where $`D`$ is the diffusion coefficient (of order the inverse mean free path), $`v_w`$ is the bubble wall velocity, and $`\mathrm{\Gamma }`$ is a damping rate representing decays or inelastic collisions of the left-handed fermions. The source term for each species is related to the classical force through a thermal average
$$S(x)=\frac{v_wD}{\stackrel{}{v}^{\mathrm{\hspace{0.17em}2}}}v_xF(x)^{}.$$
(3)
Although this looks similar to spontaneous baryogenesis , it is not the same; the latter arises through the collision term in the Boltzmann equation, while ours comes from the flow term , and we find that it gives a somewhat larger effect.
Once the left-handed quark density is found from eq. (2), the density of baryons generated by sphalerons can be computed by integrating $`\mu _L(x)_{q_i}\mu _{L}^{}{}_{i}{}^{}`$ in front of the wall:
$$n_B=\frac{9\mathrm{\Gamma }_{\mathrm{sph}}T^2}{2v_w}_0^{\mathrm{}}\mu _L(x)e^{c_b\mathrm{\Gamma }_{\mathrm{sph}}x/v_w}𝑑x,$$
(4)
where $`\mathrm{\Gamma }_{\mathrm{sph}}`$ is the diffusion rate of the Chern-Simons number as measured by lattice simulations : $`\mathrm{\Gamma }_{\mathrm{sph}}=(20\pm 2)\alpha _w^5T`$. The exponential accounts for sphaleron-induced relaxation of the baryon asymmetry back to zero due to restoration of thermal equilibrium, in the limit of a very slowly moving wall. The coefficient $`c_b`$ depends on the squark spectrum: it equals to $`45/4`$ if all squarks are heavy and $`72/7`$ if only the right handed stop is light.
3. When the above picture is adapted to the MSSM some complications occur, because the quarks do not get complex masses and hence are not directly sourced. The chargino mass matrix does contain complex phases however, and the chiral asymmetry which develops in the chargino sector induces one for the quarks because of the strong coupling between Higgsinos and the top quark. This system involves not just one, but many coupled diffusion equations. We will show that they can, nevertheless, be reduced to a single equation by reasonable approximations.
We start the treatment of the MSSM by deriving the source term analogous to $`S(x)`$ in eqs. (23). The mass term in the Lagrangian for the charginos is
$$\overline{\psi }_RM_\chi \psi _L=(\overline{\stackrel{~}{w}^^+},\overline{\stackrel{~}{h}_2^^+})_R\left(\begin{array}{cc}m_2& gH_1\\ gH_2& \mu \end{array}\right)\left(\begin{array}{c}\stackrel{~}{w}^^+\\ \stackrel{~}{h}_1^^+\end{array}\right)_L$$
(5)
plus the Hermitian conjugate, $`\overline{\psi }_LM_\chi ^{}\psi _R`$. Because there are two Higgs doublets, $`H_1`$ and $`H_2`$, there are two corresponding Higgsino fields $`\stackrel{~}{h}_{1,2}^^+`$. $`\stackrel{~}{w}^^+`$ is the wino, superpartner of the $`W`$ boson, and $`g`$ is the weak gauge coupling. The complex phases of the wino mass $`m_2`$ and the $`\mu `$ parameter are the origin of a CP violating force. By again solving the Dirac equation in the WKB approximation, one can find the spin-dependent forces which act on the two mass eigenstates, with masses $`m_\pm ^2`$. Since Higgsinos couple strongly to top quarks, we are interested in the one which smoothly evolves into the Higgsino state in front of the wall, where $`H_{1,2}0`$. The spin-dependent part of the force has the same form as in eq. (1), with the spatially varying phase given by
$$m_\pm ^2\theta _h^{}=\frac{g^2\mathrm{Im}(m_2\mu )}{(m_+^2m_{}^2)}\left(H_1H_2^{}+H_1^{}H_2\right),$$
(6)
where the Higgsino-like mass eigenvalue, $`m_{\stackrel{~}{h}}^2(x)`$, is $`m_{}^2`$ (the lighter one) if $`|\mu |^2<|m_2|^2`$, and $`m_+^2`$ otherwise. We remark that the effective WKB expansion parameter $`\theta _h^{}/E`$ remains small even when mass gap $`m_+^2m_{}^2`$ attains its minimum value; this corresponds to $`|m_2||\mu |`$, where one has parametrically $`\theta _h^{}/EgH_i/wE|\mu |1`$, because for typical wall widths $`wE\text{ }\stackrel{>}{}\text{ }20`$ .
Having specified the source term, we must now deal with the diffusion equations in which it appears. In general, these are a set of coupled equations for the two Higgsinos ($`\stackrel{~}{h}_{1,2}`$), the winos, 6 flavors of right-handed quarks, 3 generations of left-handed quark doublets, and all the superpartners which are light enough to be present at the temperature $`T`$. Following Huet and Nelson , we will make the simplifying assumption that the supergauge interactions (e.g., the coupling of winos to quarks and squarks) are in thermal equilibrium, so that particle and corresponding sparticle chemical potentials are equal to each other, species by species; in particular gaugino chemical potentials are driven to zero. We will further assume that the right-handed top squark is the only light squark. The latter must be light to get a strong enough electroweak phase transition, and it turns out that making all the others heavy is favorable to electroweak baryogenesis, as well as satisfying constraints on CP violation in the MSSM.
With these simplifications, the diffusion equation network can be reduced to four equations, for the potentials of the Higgsinos $`\stackrel{~}{h}_1`$, $`\stackrel{~}{h}_2`$ and the left-handed third generation quark doublet $`q_3`$ and the right-handed quark $`t_R`$. Ignoring interactions involving left-handed squarks, which we assume to be heavy, the important interactions coupling these species come from the potential
$$V=y\overline{q}_3H_2t_R+\mu \stackrel{~}{h}_1\stackrel{~}{h}_2+\mathrm{h}.\mathrm{c}.$$
(7)
The rate for the Yukawa interaction is $`\mathrm{\Gamma }_y`$. In addition there are helicity flipping interactions with rate $`\mathrm{\Gamma }_{hf}`$ coupling $`\stackrel{~}{h}_1`$ and $`\stackrel{~}{h}_2`$, due to the $`\mu `$ term, and in the broken phase inside the bubble ($`x<0`$), the top quark Yukawa coupling becomes the top mass, which causes helicity flips between the left- and right-handed top quarks, $`q_3`$ and $`t_R`$, with a rate of $`\mathrm{\Gamma }_m`$. Defining the diffusion operator $`𝒟_i=6(D_i_x^2+v_w_x)`$, and defining also the convenient linear combinations $`\mu _\pm =\frac{1}{2}(\mu _{\stackrel{~}{h}_1}\pm \mu _{\stackrel{~}{h}_2})`$ and $`\mu _y=\mu _+\mu _{}\mu _{t_R}+\mu _{q_3}`$, we have
$`𝒟_h\mu _++2\mathrm{\Gamma }_{hf}\mu _++\frac{1}{2}\mathrm{\Gamma }_y\mu _y`$ $`=`$ $`S_H`$ (8)
$`𝒟_h\mu _{}\frac{1}{2}\mathrm{\Gamma }_y\mu _y`$ $`=`$ $`0`$ (9)
$`\frac{1}{3}𝒟_q\mu _{q_3}+\frac{1}{2}\kappa _{\stackrel{~}{h}}\mathrm{\Gamma }_y\mu _y\stackrel{~}{\mathrm{\Gamma }}_m+2\stackrel{~}{\mathrm{\Gamma }}_{ss}`$ $`=`$ $`0`$ (10)
$`\frac{1}{2}𝒟_q\mu _{t_R}\frac{1}{2}\kappa _{\stackrel{~}{h}}\mathrm{\Gamma }_y\mu _y+\stackrel{~}{\mathrm{\Gamma }}_m\stackrel{~}{\mathrm{\Gamma }}_{ss}`$ $`=`$ $`0.`$ (11)
In deriving these equations we have assumed that all squarks but right handed decouple from thermal equilibrium. The factor $`\kappa _{\stackrel{~}{h}}\frac{1}{2}x_{\stackrel{~}{h}}^2K_2(x_{\stackrel{~}{h}})`$, where $`x_{\stackrel{~}{h}}=m_{\stackrel{~}{h}}/T`$ and $`K_2`$ is the modified Bessel function, accounts for partial decoupling of higgsinos. The top-quark helicity-flip term is defined to be $`\stackrel{~}{\mathrm{\Gamma }}_m\mathrm{\Gamma }_m\theta (x)(\mu _{t_R}\mu _{q_3})/T`$ and that for the strong sphaleron rate is $`\stackrel{~}{\mathrm{\Gamma }}_{ss}\mathrm{\Gamma }_{ss}(20\mu _{q_3}+26\mu _{t_R})/T`$ with $`\mathrm{\Gamma }_{ss}=1500\kappa _{\mathrm{sph}}\alpha _W^5T`$ .
It is noteworthy that only the linear combination $`\mu _+`$ gets directly sourced in our WKB treatment, whereas the source for the combination $`\mu _{}`$ vanishes (in this respect ref. was in error). This is in contrast to papers which treat the particle reflections from the wall quantum mechanically; these works find that both $`\mu _+`$ and $`\mu _{}`$ are sourced. Here we point out that, within the quantum reflection formalisms, it is still true that $`\mu _+`$ has a larger source than $`\mu _{}`$, because $`H_1^{}H_2+H_2^{}H_1`$ is larger than $`H_1^{}H_2H_2^{}H_1`$ . However nobody has heretofore considered the former because of the unfortunate approximation of imposing equilibrium of the Yukawa interactions $`(\mathrm{\Gamma }_y\mathrm{})`$, which forces $`\mu _+`$ to zero.
The network of diffusion equations can be integrated numerically to to obtain $`\mu _{t_R}`$ and $`\mu _{q_3}`$. From these one can compute the left-handed quark potential $`\mu _L=18\mu _{t_R}+15\mu _{q_3}`$, which must then be numerically integrated in eq. (4) to obtain the baryon asymmetry.
4. We have carried out the above procedure to study how the BAU depends on the velocity of the bubble wall $`v_w`$, the wall width $`w`$ (appearing in the Higgs field profile as $`H_0(1\mathrm{tanh}(x/w))/2`$), $`\mathrm{tan}\beta =H_2/H_1`$ (at zero temperature) and the chargino mass parameters $`\mu `$ and $`m_2`$. We take as our preferred values $`\mathrm{tan}\beta =3`$, $`w=6/T`$ , $`v_w=0.1`$ and $`\mu =m_2=100`$ GeV. We also took $`T=90`$ GeV and $`H_0=110`$ GeV. We computed the values of $`D_{\stackrel{~}{h}}`$, $`\mathrm{\Gamma }_y`$, $`\mathrm{\Gamma }_{hf}`$ and $`\mathrm{\Gamma }_m`$ from the actual Feynman diagrams and thermally averaged cross sections, and also using the results of ref. . With our preferred parameters we find $`D_{\stackrel{~}{h}}=80/T`$, $`\mathrm{\Gamma }_m=0.0012T`$, $`\mathrm{\Gamma }_{hf}=0.001x_{\stackrel{~}{h}}^2/[(1+x_{\stackrel{~}{h}})(2+x_{\stackrel{~}{h}})]`$ and $`\mathrm{\Gamma }_y=[0.012+0.022x_h(11.87/x_h^2+0.05/x_h^4)^{3/2}]T`$. In the following figures where certain of these quantities are varied, those which are not specified have the above values. It is customary to express the BAU as a scaled ratio of baryons to photons, $`\eta _{10}=10^{10}n_B/n_\gamma =7\times 10^{10}n_B/s`$, where $`s=(2\pi ^2/45)g_{}T^3`$, and we take the number of degrees of freedom in the MSSM to be $`g_{}=110`$. Current limits from big bang nucleosynthesis give $`3\text{ }\stackrel{<}{}\text{ }\eta _{10}\text{ }\stackrel{<}{}\text{ }4`$.
In figure 1, the BAU is plotted as a function of wall velocity, varying the width of the bubble wall to obtain the different curves. We have taken the CP-violating phase appearing in Im($`m_2\stackrel{~}{\mu }`$) (eq. (6)) to be maximal; to satisfy $`\eta _{10}3`$, one should rescale this phase accordingly. The efficiency of baryogenesis tends to peak for $`v_w=0.010.1`$. Interestingly, recent work on bubble expansion has suggested just such a range of small values in the MSSM.
We can also consider the dependence of the BAU on the chargino mass parameters, $`\mu `$ and $`m_2`$. Figure 2 exhibits contours of constant $`\eta _{10}`$ assuming the maximal magnitude of the CP-violating phase, $`\mathrm{arg}(m_2\mu )=\pi /2`$, and wall velocity $`v_w=0.1`$. Since we need $`\eta _{10}3`$, and the maximum value in the figure occurs around $`\eta _{10}20`$ for $`\mu m_250`$ GeV, we find that CP phase can be no smaller than 0.15, and then only if the charginos are rather light. Such large values of the CP phase could come into conflict with the most stringent present experimental constraints, which come from searches for the electric dipole moment of the <sup>199</sup>Hg atom . In order to evade these constraints, one must assume that most of the squarks or gauginos contributing to EDM loop diagrams are heavy, in the TeV range.
5. In summary, we have presented a quantitatively accurate analysis of the baryon asymmetry of the universe in the MSSM, using the classical force mechanism, which allows a consistent computation of the sources appearing in the diffusion equations. A CP-violating force acts on Higgsinos while they cross the bubble walls causing a particle-antiparticle separation in charginos, which then gets partially transformed to a chiral quark asymmetry that biases sphalerons to produce baryons. Unlike previous authors, we have not assumed that the two components of the Higgsino, $`\stackrel{~}{h}_1`$ and $`\stackrel{~}{h}_2`$, reach complete chemical equilibrium through Yukawa and helicity-flipping interactions, and we showed that a large enhancement of the BAU can result, relative to what one would get this oversimplification. Nevertheless, after carefully computing the relevant damping rates and diffusion coefficients, and accurately solving the transport equations, we find that it is difficult to generate the observed baryon asymmetry unless CP violation in the $`\mu `$ term is close to maximal. Further details of our formalism and calculations can be found in ref. .
We thank Guy Moore for very helpful correspondence, and in particular Michael Joyce and Tomislav Prokopec for collaboration on related issues. Special thanks also to Stephan Huber and Michael Schmidt for pointing out an error in our original numerical results. KK thanks CERN for hospitality during the completion of this work. |
warning/0002/hep-ph0002301.html | ar5iv | text | # Inhaltsverzeichnis
## Einleitung
Das Standardmodell der Elementarteilchenphysik ist eine äußerst erfolgreiche Theorie. Es beschreibt auf konsistente Weise die starke, schwache und elektromagnetische Wechselwirkung, und seine Vorhersagen sind experimentell mit großer Genauigkeit bestätigt worden. Nach der Entdeckung des top-Quarks steht lediglich der Nachweis des Higgs-Bosons, dessen Existenz von der Theorie gefordert wird, noch aus.
Es gibt jedoch verschiedene Anhaltspunkte, die darauf hindeuten, daß das Standardmodell keine wirklich fundamentale Theorie darstellt, sondern vielmehr als effektive Niederenergie-Näherung einer solchen den heute experimentell zugänglichen Energiebereich beschreibt. So sind die Neutrinos im Rahmen des Standardmodells masselose Teilchen, während die Evidenz für nichtverschwindende Neutrinomassen durch Experimente wie Super-Kamiokande in den letzten Jahren deutlich zugenommen hat. Auch die beobachtete Baryonasymmetrie des Universums läßt sich nicht befriedigend erklären.
Eine weitere formale Schwäche des Standardmodells besteht in der großen Zahl freier Parameter, welche durch die Theorie nicht festgelegt sind, sondern an experimentelle Ergebnisse angepaßt werden müssen. Es gibt 18 solcher Parameter, von denen allein 13 im Fermionsektor liegen; dies sind die Fermionmassen und die Quarkmischungen. Abgesehen davon liefert das Standardmodell keine wirkliche Vereinheitlichung der drei Wechselwirkungen, da die zugehörigen Kopplungskonstanten völlig unabhängig voneinander sind.
Um diese Schwachpunkte zumindest teilweise beseitigen zu können, ist eine Reihe von Versuchen gemacht worden, das Standardmodell in eine umfassendere Theorie einzubetten. Viele von ihnen liefern jedoch kaum Vorhersagen, welche in naher Zukunft experimentell verifiziert werden können. Einen in dieser Hinsicht vielversprechenden Ansatz liefern die sogenannten Grand Unified-Theorien, deren Grundidee darin besteht, die drei Kräfte des Standardmodells in einer einzigen Wechselwirkung zu vereinheitlichen. Formal geschieht das durch die Konstruktion einer Eichtheorie mit einfacher Symmetriegruppe, welche die Standardmodell-Symmetriegruppe als Untergruppe enthält. Motivation hierfür ist in erster Linie die Tatsache, daß die skalenabhängigen Kopplungen des Standardmodells sich bei sehr großen Energien von $`10^{14}10^{15}`$ GeV näherungsweise in einem Punkt treffen. Mittels spontaner Symmetriebrechung erhält man aus der Grand Unified-Theorie als effektive Näherung bei niedrigen Energien wieder das Standardmodell.
Die wohl bemerkenswerteste Vorhersage von Grand Unified-Theorien besteht in der Instabilität des Protons und des gebundenen Neutrons aufgrund von baryon- und leptonzahlverletzenden Wechselwirkungen. Einerseits sind Zerfälle von Nukleonen bisher nicht beobachtet worden, andererseits ist die relevante Wechselwirkung in Grand Unified-Theorien wegen der großen Masse der zugehörigen Eichbosonen bei niedrigen Energien auch überaus schwach. Die einfachste der Grand Unified-Theorien, welche auf der Gruppe $`SU(5)`$ beruht, ist experimentell ausgeschlossen worden, da sie unter anderem zu kurze Lebensdauern für die Nukleonen liefert. Modelle auf der Grundlage der $`SO(10)`$ wie das hier untersuchte besitzen diesen Mangel nicht.
Da sich die laufenden Standardmodell-Kopplungen nicht exakt in einem Punkt treffen, ist eine direkte Vereinheitlichung der drei Wechselwirkungen nicht möglich; es muß eine zusätzliche intermediäre Symmetrie vorhanden sein. Die Skala, bei welcher diese Symmetrie in das Standardmodell gebrochen wird, liegt üblicherweise im Bereich $`10^{10}10^{12}`$ GeV und hat somit die richtige Größenordnung, damit über den See-Saw-Mechanismus sehr kleine, aber nichtverschwindende Neutrinomassen erzeugt werden können. Weitere Vorzüge von Grand Unified-Theorien sind Beziehungen zwischen den Massenmatrizen der Quarks und der Leptonen. Infolgedessen können durch Kenntnis der Massenmatrizen der geladenen Fermionen Vorhersagen für die Neutrinomassen und -mischungen gewonnen werden. Weiterhin kann im Rahmen von $`SO(10)`$-Modellen mit intermediärer Massenskala und schweren Majorana-Neutrinos die Baryonasymmetrie des Universums erklärt werden.
Ein anderer Versuch, den Fermionsektor des Standardmodells besser zu verstehen, besteht in der Untersuchung von phänomenologisch motivierten Ansätzen für die Massenmatrizen der Fermionen. Diese Ansätze zeichnen sich durch Symmetrieanforderungen oder sogenannte Texturen, das heißt Nullen als Matrixeinträge an bestimmten Stellen, aus. Auf die zugrundeliegende Theorie jenseits des Standardmodells, welche den Ansatz realisiert, wird im allgemeinen nicht weiter eingegangen. Damit will man die Zahl der freien Parameter des Standardmodells reduzieren und Beziehungen zwischen Massen und Mischungen erhalten.
Schließlich kann man beide Zugänge kombinieren und Massenmodelle auf der Grundlage von Grand Unified-Theorien konstruieren. Das eröffnet die Möglichkeit, die Vorzüge dieser Theorien gegenüber dem Standardmodell auszunutzen und die Ansätze wegen der gegebenen Beziehungen zwischen den Fermionmassenmatrizen weniger willkürlich zu machen. Ferner haben in solchen Modellen alle Fermionmischungen Einfluß auf zumindest prinzipiell observable Größen wie Nukleonzerfallsraten, während im Standardmodell lediglich eine bestimmte Kombination der linkshändigen Quarkmischungen, die CKM-Matrix, beobachtbar ist. Daraus folgen überprüfbare Konsequenzen, an denen man den Erfolg des Ansatzes messen kann.
Gegenstand dieser Arbeit wird ein Massenmodell im Rahmen einer $`SO(10)`$-Theorie sein. Es wird ein asymmetrischer „Nearest Neighbour Interaction“-Ansatz für die Massenmatrizen der Fermionen benutzt, welcher durch eine globale $`U(1)`$-Familiensymmetrie realisiert wird. Dieser Ansatz führt auf voneinander unabhängige rechts- und linkshändige Mischungen und bietet ausdrücklich die Möglichkeit, daß diese betragsmäßig groß sind. In vielen Massenmodellen werden große Mischungen im Bereich der geladenen Fermionen mit dem Hinweis auf die relativ kleinen CKM-Mischungen der Quarks außer Acht gelassen. In der Tat besitzen alle gefundenen Lösungen des untersuchten Modells mehrere große Mischungen, was zu Verzweigungsraten der Nukleonen führt, die sich von denen im Fall verschwindender Mischungen deutlich unterscheiden. Zusätzlich sind die erhaltenen Neutrinoeigenschaften, welche ebenfalls Modellvorhersagen darstellen, in der Lage, die Anomalien der Sonnen- und atmosphärischen Neutrinos durch Oszillationslösungen zu erklären.
In den letzten Jahren hat sich die Forschungstätigkeit hauptsächlich auf supersymmetrische Grand Unified-Theorien, welche spontan in das minimale supersymmetrische Standardmodell gebrochen werden, beschränkt. Wird die Supersymmetrie, eine Symmetrie zwischen Fermionen und Bosonen, bei vergleichsweise niedrigen Energien von $`1`$ TeV effektiv gebrochen, so treffen sich die Kopplungskonstanten des supersymmetrische Standardmodells bei etwa $`10^{16}`$ GeV genau in einem Punkt. Da die Massen der Superpartner dann nicht sehr viel größer als die der gewöhnlichen Teilchen sind, bietet sich die Möglichkeit, das Divergenzverhalten der Theorie zu verbessern und das Hierarchieproblem zu lösen. Allerdings läßt die Abwesenheit von experimentellen Hinweisen auf eine bei kleinen Energien gebrochene Supersymmetrie in der Natur ein solches Szenario zunehmend unwahrscheinlicher erscheinen. Auch eine für die Erzeugung von Neutrinomassen über den See-Saw-Mechanismus erforderliche intermediäre Massenskala läßt sich in supersymmetrischen Grand Unified-Theorien nicht auf natürliche Weise realisieren. Desweiteren sind supersymmetrische Modelle in ihrer Vorhersagekraft durch den Einfluß zahlreicher unbekannter Größen, wie zum Beispiel die Massen der Superpartner, stark eingeschränkt. Deshalb wird hier der Standpunkt vertreten, daß Modelle ohne Supersymmetrie weiter untersucht werden sollten. Dies schließt eine Brechung der Supersymmetrie bei sehr hohen Energien keineswegs aus.
Im ersten Kapitel wird zunächst ein Überblick über die Grundkonzepte und wichtigsten Eigenschaften des Standardmodells gegeben, auch dessen Grenzen werden genauer betrachtet. Kapitel 2 behandelt den Themenkomplex der Grand Unified-Theorien. Nach einer Schilderung der grundlegenden Ideen und der allgemeinen Vorgehensweise bei der Konstruktion solcher Modelle werden die $`SU(5)`$ als Prototyp der Grand Unified-Theorien und die $`SO(10)`$ als das dieser Arbeit zugrundeliegende Beispiel ausführlich diskutiert; desweiteren wird auf Vor- und Nachteile der Modelle eingegangen. Gegenstand des dritten Kapitels sind die theoretischen Grundlagen und der experimentelle Status von Neutrino-Oszillationen. Die drei bis heute beobachteten Neutrino-Anomalien und ihre möglichen Oszillationslösungen durch massive Neutrinos werden vorgestellt. Kapitel 4 beginnt mit vorbereitenden Arbeiten wie der Bestimmung der Symmetriebrechungsskalen und Kopplungen. Mittelpunkt des Kapitels ist der Ansatz für das betrachtete $`SO(10)`$-Massenmodell und dessen numerische Lösung. Als ein wichtiges Resultat erhält man daraus Voraussagen über die Eigenschaften im Neutrinosektor der Theorie, welche sich als phänomenologisch sinnvoll erweisen. In Kapitel 5 schließlich werden für die analysierten Lösungen die partiellen und totalen Zerfallsraten der Nukleonen berechnet. Diese stellen eine wesentliche und in absehbarer Zeit experimentell überprüfbare Vorhersage des untersuchten Modells dar.
Ziel und Motivation dieser Untersuchung lassen sich abschließend mit den einleitenden Worten aus sehr treffend zusammenfassen:
„The Standard Model is unlikely to be a fundamental theory; it contains 18 arbitrary parameters, 13 of which are the fermion masses and mixing angles. In a fundamental theory, these should be calculable from a few inputs just as the hydrogen spectral lines follow from Quantum Mechanics. We are very far from such a theory of fermion masses. We would be fortunate to have an analogue of Balmer’s formula since it might lead us to the fundamental theory. The framework described here is, at best, an attempt to obtain such a formula.“
## Kapitel 1 Das Standardmodell
### 1.1 Grundkonzepte des Standardmodells
Das Standardmodell (SM) der Elementarteilchenphysik ist eine renormierbare Eichtheorie, welche die Theorie der starken Wechselwirkung, die Quantenchromodynamik (QCD), und das Glashow-Weinberg-Salam-Modell der elektroschwachen Wechselwirkung zusammenfaßt. Es basiert auf der Invarianz unter lokalen $`SU(3)_CSU(2)_LU(1)_Y`$-Eichtransformationen; die Symmetriegruppe $`G_{\text{SM}}`$ ist ein direktes Produkt aus drei Faktoren.
Die QCD beruht auf der Eichgruppe $`SU(3)_C`$ und beschreibt die Wechselwirkung von Teilchen mit Farbladung, den Quarks und Gluonen. Die Quarks sind die fermionischen Grundbausteine der stark wechselwirkenden Materie und treten in drei verschiedenen Farbzuständen auf. Unter $`SU(3)_C`$-Transformationen verhalten sie sich wie die Fundamentaldarstellung $`\mathrm{𝟑}`$. In der Natur werden allerdings nur farbneutrale Kombinationen von drei Quarks (Baryonen) bzw. Quark und Antiquark (Mesonen) beobachtet; dieses Phänomen wird als Confinement bezeichnet. Die Gluonen sind die mit den acht Generatoren der $`SU(3)_C`$ verbundenen Vektorbosonen und transformieren sich gemäß der adjungierten Darstellung. Durch den Austausch von Gluonen wird die starke Wechselwirkung zwischen den Quarks vermittelt.
Die Theorie der elektroschwachen Wechselwirkung besitzt eine $`SU(2)_LU(1)_Y`$-Eichsymmetrie, die durch den Higgs-Mechanismus spontan in die $`U(1)_{\text{em}}`$ der Quantenelektrodynamik (QED) gebrochen ist. Eine Grundeigenschaft dieses Modells ist die Paritätsverletzung, da das Transformationsverhalten der Fermionen von deren Chiralität abhängt. Man zerlegt die durch Dirac-Spinoren $`\mathrm{\Psi }`$ dargestellten Fermionfelder gemäß
$$\mathrm{\Psi }_{L,R}=\frac{1}{2}(1\gamma _5)\mathrm{\Psi },\mathrm{\Psi }=\mathrm{\Psi }_L+\mathrm{\Psi }_R$$
(1.1)
in ihre links- und rechtshändigen Komponenten, sogenannte Weyl-Spinoren. Die linkshändigen Quarks und Leptonen sind in $`SU(2)_L`$-Doubletts angeordnet, während die rechtshändigen Fermionen Singuletts bilden und an der schwachen Wechselwirkung nicht teilnehmen.
Die elektrische Ladung der Teilchen ergibt sich aus ihrer Hyperladung $`Y`$ und der Komponente $`T_3`$ des schwachen Isospins ($`T_a=\sigma _a/2`$ mit $`a=\mathrm{1,2,3}`$ sind die Generatoren der $`SU(2)_L`$-Lie-Algebra) aus der Beziehung
$$Q=T_3+Y$$
(1.2)
In Erweiterungen des SM erweist es sich aus gruppentheoretischen Gründen als sinnvoll, statt der rechtshändigen Teilchen die linkhändigen Komponenten der Antiteilchen zu betrachten. Die Antiteilchen erhält man durch die Anwendung der Ladungskonjugation $`𝒞`$
$$𝒞\mathrm{\Psi }𝒞^1\mathrm{\Psi }^C=C\overline{\mathrm{\Psi }}^T\text{mit}C=i\gamma _2\gamma _0,\mathrm{\Psi }_{L,R}^C(\mathrm{\Psi }_{R,L})^C$$
(1.3)
Weiterhin gelten die Identitäten
$$\mathrm{\Psi }_{R,L}=C(\overline{\mathrm{\Psi }}_{L,R}^C)^T\text{und}\overline{\mathrm{\Psi }}_{R,L}=(\mathrm{\Psi }_{L,R}^C)^TC$$
(1.4)
Die Fermionen treten im SM in drei Familien mit jeweils gleichen Quantenzahlen auf. Jede dieser Familien transformiert sich nach der Darstellung
$$(\mathbf{3,2})_{{\scriptscriptstyle \frac{1}{6}}}(\overline{\mathrm{𝟑}}\mathbf{,1})_{{\scriptscriptstyle \frac{2}{3}}}(\overline{\mathrm{𝟑}}\mathbf{,1})_{{\scriptscriptstyle \frac{1}{3}}}(\mathbf{1,2})_{{\scriptscriptstyle \frac{1}{2}}}(\mathbf{1,1})_1$$
(1.5)
der SM-Symmetriegruppe $`G_{\text{SM}}`$. Der erste Eintrag bezeichnet hierbei die $`SU(3)_C`$-Darstellung, der zweite die $`SU(2)_L`$-Darstellung und der Index die $`U(1)_Y`$ -Hyperladungsquantenzahl (die komplex konjugierte Darstellung wird durch einen Querstrich gekennzeichnet). In Tabelle 1.1 ist der fermionische Teilcheninhalt des SM aufgeführt.
Es ist zu beachten, daß im Fermionspektrum des SM keine rechtshändigen Neutrinos vorkommen, da diese sich aufgrund ihrer Farb- und Ladungsneutralität nach der $`G_{\text{SM}}`$-Darstellung $`(\mathbf{1,1})_0`$ transformieren würden, also an keiner SM-Wechselwirkung teilnähmen. Dies steht im Einklang mit der Tatsache, daß rechtshändige Neutrinos und linkshändige Antineutrinos in Experimenten nicht beobachtet werden.
Die Eichbosonen, die als Austauschteilchen die Wechselwirkung vermitteln, transformieren sich stets wie die adjungierte Darstellung der Symmetriegruppe, im Falle des SM also gemäß $`(\mathbf{8,1})_0(\mathbf{1,3})_0(\mathbf{1,1})_0`$. Die Kopplungsstärken der SM-Eichbosonen an die fermionischen Ströme $`j^\mu =\overline{\mathrm{\Psi }}\gamma ^\mu \mathrm{\Psi }`$ werden mit $`g_3`$, $`g_2`$ und $`g^{}`$ bezeichnet. Ferner besitzen Eichbosonen in Modellen mit nichtabelscher Symmetriegruppe auch eine Selbstwechselwirkung. In Erweiterungen des SM durch Grand Unified-Theorien (GUTs) wird statt $`g^{}`$ im allgemeinen mit $`g_1=\sqrt{5/3}g^{}`$ gearbeitet, da $`g_1`$ die korrekte Normierung besitzt. Die Eigenschaften der SM-Eichbosonen sind in Tabelle 1.2 zusammengefaßt.
### 1.2 Renormierung und laufende Kopplungen
Ein wichtiges Kriterium für die formale Konsistenz einer Eichtheorie ist ihre Renormierbarkeit. Dies bedeutet, daß bei der Berechnung von physikalischen Prozessen in höheren Ordnungen der Störungstheorie nur endlich viele qualitativ verschiedene Divergenzen auftreten. Diese können dann durch eine Redefinition der Modellparameter wie Kopplungen und Massen absorbiert werden .
Die Unendlichkeiten treten in Form von divergenten Impulsintegralen auf. Durch Anwendung eines Regularisierungsverfahrens werden die Integrale durch Ausdrücke ersetzt, die von einem neuen Parameter, dem Regularisierungsparameter, abhängen, und für einen bestimmten Wert desselben wieder die ursprüngliche Gestalt annehmen. Bei der dimensionalen Regularisierung zum Beispiel werden die Integrale statt in vier in $`D=42\epsilon `$ Impulsraumdimensionen gelöst; die Divergenzen gehen dann in Terme $`1/\epsilon `$ über. Im anschließenden Renormierungsprozeß werden die divergenten Parameter $`𝒢_0`$ in der Lagrangedichte durch Einführung von Renormierungskonstanten $`Z`$, welche die $`1/\epsilon `$-Terme aufnehmen, in die endlichen renormierten Größen $`𝒢_R`$ umgewandelt: $`𝒢_0Z𝒢_R`$.
Im Rahmen der Regularisierung wird aus formalen Gründen zwangsläufig eine freie Massenskala $`\mu `$ in die Theorie eingeführt. Sowohl die Renormierungskonstanten als auch die renormierten Größen hängen von dieser Skala ab; man spricht von laufenden Größen. Die funktionalen Zusammenhänge, welche die Skalenabhängigkeit der renormierten Größen beschreiben, werden als Renormierungsgruppengleichungen bezeichnet . Sie sind immer dann von Bedeutung, wenn physikalische Größen bei verschiedenen Massen- bzw. Energieskalen miteinander verglichen werden. Die in dieser Arbeit verwendeten Renormierungsgruppengleichungen sind in Anhang B angegeben. Die Energieabhängigkeit observabler Größen ist experimentell bestätigt; so hat die Kopplungsstärke $`\alpha _{\text{em}}=e^2/4\pi `$ der QED bei $`\mu 0`$ den Betrag 1/137, bei $`\mu M_Z`$ ist sie etwa 1/129 groß .
Die Renormierbarkeit einer Eichtheorie kann durch das Auftreten von Anomalien zerstört werden. Anomalien sind Symmetrien der klassischen Lagrangedichte, die durch den Prozeß der Quantisierung gebrochen werden . In Modellen mit chiralen Fermionen äußert sich dies durch das Auftreten von linear divergenten Strahlungskorrekturen in Form von fermionischen Dreiecksdiagrammen mit einer ungeraden Anzahl axialer Vertizes $`\gamma _\mu \gamma _5`$. Die Divergenzen verletzen die Slavnov-Taylor-Identitäten, deren Gültigkeit für die vollständige Renormierung der Theorie in allen Ordnungen der Störungsrechnung benötigt wird. Diese Beiträge sind im wesentlichen proportional zu
$$A_{abc}=A_{abc}^LA_{abc}^R=\text{Tr}\left(\{T_a,T_b\}T_c\right)_L\text{Tr}\left(\{T_a,T_b\}T_c\right)_R,$$
(1.6)
wobei die $`T_i`$ die Generatoren der zur Symmetriegruppe gehörenden Lie-Algebra sind, und zwar in der Darstellung, nach der sich die links- beziehungsweise rechtshändigen Fermionen transformieren . Die Renormierbarkeit ist demnach sichergestellt, wenn
$$A_{abc}^{L,R}=\text{Tr}\left(\{T_a,T_b\}T_c\right)_{L,R}\stackrel{!}{=}0$$
(1.7)
gilt. Im Rahmen des SM führt (1.7) auf die vier Bedingungen
$`\text{Tr}\left(\{T_a,T_b\}T_c\right)`$ $`=`$ $`0,\text{Tr}\left(T_aY^2\right)=\mathrm{\hspace{0.33em}0},`$ (1.8)
$`\text{Tr}\left(\{T_a,T_b\}Y\right)`$ $`=`$ $`0,\text{Tr}\left(Y^3\right)=\mathrm{\hspace{0.33em}0}`$
Die beiden oberen Gleichungen sind aufgrund der Spurfreiheit der $`SU(2)_L`$-Generatoren automatisch erfüllt, die unteren beiden wegen der SM-Zuordnung der Hyperladungen in (1.5). Das SM ist also deshalb anomaliefrei, weil sich die Beiträge der Quarks und der Leptonen gerade aufheben.
### 1.3 Symmetriebrechung und Massenerzeugung
In Eichtheorien können prinzipiell zwei Arten von Massentermen vorkommen. Wenn $`\mathrm{\Psi }`$ ein Dirac-Spinor ist und $`\mathrm{\Psi }_{L,R}`$ seine links- und rechtshändigen Komponenten bezeichnet, so hat der Dirac-Massenterm die Gestalt (h.c. steht für hermitesch konjugiert)
$$m\overline{\mathrm{\Psi }}\mathrm{\Psi }=m(\overline{\mathrm{\Psi }}_L\mathrm{\Psi }_R+\overline{\mathrm{\Psi }}_R\mathrm{\Psi }_L)=m(\overline{\mathrm{\Psi }}_L\mathrm{\Psi }_R+\text{h.c.})$$
(1.9)
Das ist derselbe Massenterm, der auch in der Dirac-Gleichung vorkommt. Unter Verwendung von (1.4) kann man $`\overline{\mathrm{\Psi }}_R\mathrm{\Psi }_L`$ auch als $`(\mathrm{\Psi }_L^C)^TC\mathrm{\Psi }_L`$ schreiben.
Desweiteren existieren die lorentzinvarianten Größen
$`{\displaystyle \frac{M_R}{2}}(\overline{\mathrm{\Psi }}_L^C\mathrm{\Psi }_R+\text{h.c.})`$ $`=`$ $`{\displaystyle \frac{M_R}{2}}(\overline{\mathrm{\Psi }}_L^CC(\overline{\mathrm{\Psi }}_L^C)^T+\text{h.c.})\text{und}`$ (1.10)
$`{\displaystyle \frac{M_L}{2}}(\overline{\mathrm{\Psi }}_L\mathrm{\Psi }_R^C+\text{h.c.})`$ $`=`$ $`{\displaystyle \frac{M_L}{2}}(\overline{\mathrm{\Psi }}_LC(\overline{\mathrm{\Psi }}_L)^T+\text{h.c.}),`$ (1.11)
die man als Majorana-Massenterme bezeichnet, da sie nur für Majorana-Teilchen definiert sind. Letztere werden durch Spinoren $`\mathrm{\Psi }`$ beschrieben, für welche $`\mathrm{\Psi }^C\mathrm{\Psi }`$ gilt, das heißt Teilchen und Antiteilchen sind identisch. Deswegen sind Majorana-Teilchen zwangsläufig elektrisch neutral. Majorana-Massenterme werden später im Zusammenhang mit Neutrinomassen von Bedeutung sein.
Das Transformationsverhalten von Massentermen für die geladenen Fermionen im SM sieht folgendermaßen aus:
$`\text{(}u,c,t\text{)-Quarks:}(\mathbf{3,2})_{{\scriptscriptstyle \frac{1}{6}}}(\overline{\mathrm{𝟑}}\mathbf{,1})_{{\scriptscriptstyle \frac{2}{3}}}`$ $`=`$ $`(\mathbf{1,2})_{{\scriptscriptstyle \frac{1}{2}}}(\mathbf{8,2})_{{\scriptscriptstyle \frac{1}{2}}}`$ (1.12)
$`\text{(}d,s,b\text{)-Quarks:}(\mathbf{3,2})_{{\scriptscriptstyle \frac{1}{6}}}(\overline{\mathrm{𝟑}}\mathbf{,1})_{{\scriptscriptstyle \frac{1}{3}}}`$ $`=`$ $`(\mathbf{1,2})_{{\scriptscriptstyle \frac{1}{2}}}(\mathbf{8,2})_{{\scriptscriptstyle \frac{1}{2}}}`$ (1.13)
$`\text{geladene Leptonen:}(\mathbf{1,2})_{{\scriptscriptstyle \frac{1}{2}}}(\mathbf{1,1})_1`$ $`=`$ $`(\mathbf{1,2})_{{\scriptscriptstyle \frac{1}{2}}}`$ (1.14)
Die Lagrangedichte des SM kann demnach keine Fermionmassenterme enthalten, da diese nicht invariant unter $`G_{\text{SM}}`$-Transformationen sind. Das steht aber im Widerspruch zu der experimentellen Beobachtung massiver Teilchen.
Abgesehen davon wird in der Natur keineswegs die Symmetrie des SM, sondern eine $`SU(3)_CU(1)_{\text{em}}`$-Symmetrie beobachtet.
Beide Probleme können durch die spontane Brechung lokaler Eichsymmetrien, den sogenannten Higgs-Mechanismus, gelöst werden . Dazu werden dem Teilchenspektrum der Eichtheorie Spin-0-Teilchen hinzugefügt. Im SM wird ein Doublett
$$\mathrm{\Phi }=\left(\begin{array}{c}\varphi ^+\\ \varphi ^0\end{array}\right)$$
(1.15)
dieser Higgs-Bosonen eingeführt, wobei $`\varphi ^+`$ und $`\varphi ^0`$ komplexe Skalarfelder sind; sie liegen in der $`G_{\text{SM}}`$-Darstellung $`(\mathbf{1,2})_{{\scriptscriptstyle \frac{1}{2}}}`$. Das Potential für $`\mathrm{\Phi }`$ hat die Form
$$V(\mathrm{\Phi })=m^2\mathrm{\Phi }^{}\mathrm{\Phi }+\lambda (\mathrm{\Phi }^{}\mathrm{\Phi })^2$$
(1.16)
und erzeugt für $`m^2<0`$ einen nichtverschwindenden Vakuumerwartungswert
$$0|\mathrm{\Phi }|0=\left(\begin{array}{c}0\\ \upsilon /\sqrt{2}\end{array}\right)\text{mit}\upsilon =\sqrt{\frac{m^2}{\lambda }}$$
(1.17)
von $`\mathrm{\Phi }`$. Den Wert von $`\upsilon `$ kann man in niedrigster Ordnung aus der meßbaren Fermi-Konstanten $`G_F`$ und der Beziehung $`\upsilon =(\sqrt{2}G_F)^{1/2}`$ bestimmen; er beträgt $`\upsilon =246.22`$ GeV. Nun ist der Vakuumzustand nicht mehr $`SU(2)_LU(1)_Y`$-symmetrisch; die Symmetrie der elektroschwachen Wechselwirkung wird spontan in die $`U(1)_{\text{em}}`$ der QED gebrochen:
$$SU(3)_CSU(2)_LU(1)_Y\stackrel{\mathrm{\Phi }}{}SU(3)_CU(1)_{\text{em}}$$
(1.18)
Durch ihre Wechselwirkung mit den Higgs-Teilchen erhalten die Eichbosonen Massenterme. Aus den vier $`SU(2)_LU(1)_Y`$-Bosonen entstehen so die physikalischen Eichbosonen, das heißt die Masseneigenzustände
$`W^\pm `$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}(W_1\pm iW_2)`$ (1.19)
$`Z`$ $`=`$ $`\mathrm{sin}\theta _WB+\mathrm{cos}\theta _WW_3`$ (1.20)
$`A`$ $`=`$ $`\mathrm{cos}\theta _WB\mathrm{sin}\theta _WW_3`$ (1.21)
$`\theta _W`$ ist der Weinberg-Winkel; für ihn gilt
$$\mathrm{sin}\theta _W=\frac{g^{}}{\sqrt{g_{}^{}{}_{}{}^{2}+g_2^2}},\mathrm{tan}\theta _W=\frac{g^{}}{g_2}$$
(1.22)
$`A`$ ist das Eichboson der $`U(1)_{\text{em}}`$, das masselose Photon. Seine Kopplungskonstante, die Elementarladung $`e`$, ergibt sich aus $`e=g_2\mathrm{sin}\theta _W`$. Die $`W`$\- und $`Z`$-Bosonen erhalten die Massen
$$M_W=\frac{g_2\upsilon }{2},M_Z=\frac{M_W}{\mathrm{cos}\theta _W},$$
(1.23)
während die Gluonen masselos bleiben, da die Higgs-Teilchen farbneutral sind. Tabelle 1.3 enthält die Werte für die Massen und den Weinberg-Winkel:
Die experimentell meßbaren Eichkopplungen der $`SU(3)_CU(1)_{\text{em}}`$-Theorie haben bei der Skala $`\mu =M_Z`$ folgende Werte ($`\alpha _3=g_3^2/4\pi `$ und $`\alpha _{\text{em}}=e^2/4\pi `$):
Daraus kann man mit Hilfe der Beziehungen
$$\alpha _1(M_Z)=\frac{5\alpha _{\text{em}}(M_Z)}{3\mathrm{cos}^2\theta _W(M_Z)},\alpha _2(M_Z)=\frac{\alpha _{\text{em}}(M_Z)}{\mathrm{sin}^2\theta _W(M_Z)},$$
(1.24)
die Werte der SM-Größen $`\alpha _{\mathrm{1,2}}(M_Z)`$ berechnen:
Von den vier reellen Freiheitsgraden in $`\mathrm{\Phi }`$ ist nur einer physikalisch, er gehört zum elektrisch neutralen Higgs-Boson $`H`$. Es ist das einzige SM-Teilchen, das experimentell noch nicht nachgewiesen wurde. Für seine Masse $`m_H=\sqrt{\lambda }\upsilon `$ gibt es lediglich eine Untergrenze: $`m_H>89.7`$ GeV .
Die Fermionmassen werden durch Einführung von Yukawa-Kopplungen zwischen den Fermionen und den Higgs-Teilchen realisiert. Diese Yukawa-Terme können wegen (1.12)-(1.14) und $`\mathrm{\Phi }(\mathbf{1,2})_{{\scriptscriptstyle \frac{1}{2}}}`$ eichinvariant konstruiert werden und haben die Form
$$_M=\underset{a,b=1}{\overset{3}{}}\left(\overline{Q}_L^a\stackrel{~}{\mathrm{\Phi }}(𝐘_u)_{ab}u_R^b+\overline{Q}_L^a\mathrm{\Phi }(𝐘_d)_{ab}d_R^b+\overline{L}_L^a\mathrm{\Phi }(𝐘_e)_{ab}e_R^b\right)+\text{h.c.}$$
(1.25)
mit
$$Q_L^a=\left(\begin{array}{c}u_L^a\\ d_L^a\end{array}\right),L_L^a=\left(\begin{array}{c}\nu _L^a\\ e_L^a\end{array}\right),\stackrel{~}{\mathrm{\Phi }}=i\tau _2\mathrm{\Phi }^{}=\left(\begin{array}{c}\varphi ^0\\ \varphi ^{}\end{array}\right)(\mathbf{1,2})_{{\scriptscriptstyle \frac{1}{2}}}$$
(1.26)
Die Indizes $`a`$ und $`b`$ bezeichnen die Fermionfamilien ($`d_L^as_L`$ für $`a=2`$), die Elemente der (3$`\times `$3)-Matrizen $`𝐘_i`$ sind die Yukawa-Kopplungen. Die Fermionmassen entstehen im Rahmen der spontanen Symmetriebrechung, wenn $`\mathrm{\Phi }`$ seinen Vakuumerwartungswert $`\upsilon `$ ausbildet. Dann sind die Massenmatrizen durch
$$𝐌_u=\frac{1}{\sqrt{2}}\upsilon 𝐘_u,𝐌_d=\frac{1}{\sqrt{2}}\upsilon 𝐘_d,𝐌_e=\frac{1}{\sqrt{2}}\upsilon 𝐘_e$$
(1.27)
gegeben; (1.25) geht über in
$$_M=\underset{a,b=1}{\overset{3}{}}\left((𝐌_u)_{ab}\overline{u}_L^au_R^b+(𝐌_d)_{ab}\overline{d}_L^ad_R^b+(𝐌_e)_{ab}\overline{e}_L^ae_R^b\right)+\text{h.c.}$$
(1.28)
Da die Massenmatrizen beliebige komplexe (3$`\times `$3)-Matrizen sein können, sind die Fermionen in Tabelle 1.1, die Eigenzustände der $`SU(2)_LU(1)_Y`$-Wechselwirkung, im allgemeinen nicht mehr mit den physikalischen Teilchen definierter Masse identisch. Letztere erhält man, wenn man die Massenmatrizen durch biunitäre Transformationen diagonalisiert:
$$𝐔_L^{}𝐌_u𝐔_R=𝐌_u^{\left(D\right)},𝐃_L^{}𝐌_d𝐃_R=𝐌_d^{\left(D\right)},𝐄_L^{}𝐌_e𝐄_R=𝐌_e^{\left(D\right)}$$
(1.29)
Die nichtverschwindenden Elemente der Diagonalmatrizen $`𝐌_i^{\left(D\right)}`$ sind die Fermionmassen; ihre Werte sind in Tabelle 1.6 aufgelistet.
Die Yukawa-Kopplungen sind ebenso wie die Eichkopplungen renormierte Größen und hängen somit von der Massenskala $`\mu `$ ab. Über (1.27) und (1.29) sind demnach auch die Fermionmassen und -mischungen skalenabhängig.
Die Transformationen (1.29) liefern den Zusammenhang zwischen den Eigenzuständen der SM-Wechselwirkungen, hier mit dem Index $`(0)`$ bezeichnet, und den Masseneigenzuständen:
$$u_{L,R}^{a\left(0\right)}=𝐔_{L,R}u_{L,R}^a,d_{L,R}^{a\left(0\right)}=𝐃_{L,R}d_{L,R}^a,e_{L,R}^{a\left(0\right)}=𝐄_{L,R}e_{L,R}^a$$
(1.30)
Nun ist (1.30) auch in den übrigen Termen der Langrangedichte auszuführen, in denen Fermionfelder vorkommen. Während die neutralen Ströme, die an das $`Z`$-Boson und das Photon koppeln, invariant unter (1.30) sind, ändern sich die geladenen schwachen Ströme:
$`\overline{u}_L^{a\left(0\right)}\gamma ^\mu d_L^{a\left(0\right)}W_\mu ^+`$ $``$ $`\overline{u}_L^a\gamma ^\mu \left(𝐔_L^{}𝐃_L\right)_{ab}d_L^bW_\mu ^+`$ (1.31)
$`\overline{d}_L^{a\left(0\right)}\gamma ^\mu u_L^{a\left(0\right)}W_\mu ^{}`$ $``$ $`\left(𝐃_L^{}𝐔_L\right)_{ab}\overline{d}_L^b\gamma ^\mu u_L^aW_\mu ^{}`$ (1.32)
Im Experiment tritt nur die (unitäre) Kombination $`𝐕𝐔_L^{}𝐃_L`$ der Mischungsmatrizen in Erscheinung; man nennt sie Cabibbo-Kobayashi-Maskawa-Matrix (CKM-Matrix) . Alle anderen Mischungen, insbesondere die rechtshändigen, sind im SM nicht observabel. Die experimentellen Grenzen für die Beträge der Elemente von $`𝐕`$ liegen bei :
$$|𝐕|=\left(\begin{array}{ccc}0.97450.9760\hfill & 0.2170.224\hfill & 0.00180.0045\hfill \\ 0.2170.224\hfill & 0.97370.9753\hfill & 0.0360.042\hfill \\ 0.0040.013\hfill & 0.0350.042\hfill & 0.99910.9994\hfill \end{array}\right)$$
(1.33)
Die in dieser Arbeit verwendete Parametrisierung für $`𝐕`$ lautet:
$$𝐕=\left(\begin{array}{ccc}C_{12}C_{31}& S_{12}C_{31}& S_{31}e^{i\delta }\\ S_{12}C_{23}C_{12}S_{23}S_{31}e^{i\delta }& C_{12}C_{23}S_{12}S_{23}S_{31}e^{i\delta }& S_{23}C_{31}\\ S_{12}S_{23}C_{12}C_{23}S_{31}e^{i\delta }& C_{12}S_{23}S_{12}C_{23}S_{31}e^{i\delta }& C_{23}C_{31}\end{array}\right)$$
(1.34)
mit $`C_{ij}=\mathrm{cos}\theta _{ij}`$ und $`S_{ij}=\mathrm{sin}\theta _{ij}`$. Als numerische Werte für die Winkel werden $`\theta _{12}=0.223`$, $`\theta _{23}=0.039`$ und $`\theta _{31}=0.003`$ benutzt. Ferner wird $`\delta =0`$ gewählt, da ein Modell mit reellen Massenmatrizen (und somit orthogonalen Mischungsmatrizen) Gegenstand der Untersuchung ist; auf das Problem der $`CP`$-Verletzung soll hier nicht näher eingegangen werden. Die Parametrisierung (1.34) mit $`\delta =0`$ wird auch für alle Mischungsmatrizen in (1.29) benutzt.
Die leptonischen Anteile der geladenen schwachen Ströme ändern sich im Falle masseloser Neutrinos nicht, da die linkshändigen Neutrinos derselben Transformation $`𝐄_L`$ wie die geladenen Leptonen unterzogen werden können.
### 1.4 Grenzen des Standardmodells
Das SM ist nicht nur mathematisch konsistent, sondern auch phänomenologisch überaus erfolgreich. Viele seiner Vorhersagen sind auf eindrucksvolle Weise experimentell bestätigt worden, zum Beispiel die Existenz der dritten Fermiongeneration, der massiven $`W^\pm `$\- und $`Z`$-Bosonen und des neutralen schwachen Stroms. Lediglich die beobachtete Baryonasymmetrie des Universums und die durch die Neutrinoexperimente der letzten Jahre implizierte Existenz massiver Neutrinos lassen sich im Rahmen des SM nicht befriedigend erklären. Davon abgesehen ist die Übereinstimmung der theoretischen Vorhersagen des SM mit den experimentellen Ergebnissen sehr gut .
Es gibt allerdings auch eine Reihe von ungeklärten theoretischen Fragestellungen, die stark darauf hindeuten, daß das SM keine wirklich fundamentale Theorie ist, sondern lediglich die effektive Näherung einer solchen für niedrige Energien.
Zunächst ist die Zahl der freien Parameter im SM, deren Werte von der Theorie nicht vorhergesagt werden, sondern an experimentelle Resultate angepaßt werden müssen, sehr groß:
* neun Massen der geladenen Fermionen
* drei Winkel und eine Phase in der CKM-Matrix
* drei Eichkopplungen
* die Higgs-Masse und -kopplungskonstante
* zwei $`\theta `$-Parameter in den $`CP`$-verletzenden Lagrangedichte-Termen $`\theta \text{Tr}F_{\mu \nu }\stackrel{~}{F}^{\mu \nu }`$
Unter Berücksichtigung massiver Neutrinos kommen drei Neutrinomassen sowie drei Winkel und drei Phasen in der leptonischen Mischungsmatrix hinzu. Eine grundlegende Theorie sollte dagegen mit möglichst wenigen freien Parametern auskommen.
Weitere Aspekte, die auf eine Theorie jenseits des SM schließen lassen, sind:
* Die Symmetriegruppe $`G_{\text{SM}}`$ ist ein direktes Produkt, woraus die Existenz dreier voneinander unabhängiger und betragsmäßig stark unterschiedlicher Kopplungskonstanten folgt; eine Vereinheitlichung der Wechselwirkungen erfolgt nicht. Ebenfalls damit verbunden ist die komplizierte Darstellung (1.5), in der die Fermionen einer Familie liegen.
* Die Zuordnung der Hyperladungen zu den Fermionen ist weitgehend willkürlich und lediglich durch (1.2) und (1.8) eingeschränkt. Dies liefert aber keine Erklärung für die Quantisierung der elektrischen Ladung in Einheiten von $`e/3`$.
* Die bis heute beobachteten Fermionen liegen in drei Familien, die hinsichtlich ihrer Quantenzahlen und ihrer Wechselwirkungen völlig identisch sind, sich aber bezüglich ihrer Massen beträchtlich unterscheiden.
* Der Hauptgrund für die Einordnung der linkshändigen Fermionen in $`SU(2)_L`$-Doubletts und der rechtshändigen in Singuletts liegt im phänomenologischen Erfolg dieses Ansatzes.
* Neutrinomassen können im Rahmen des SM zwar durch Einführung von nichtwechselwirkenden rechtshändigen Neutrinos konstruiert werden, aber es stellt sich dann die Frage, warum ihre Massen sehr viel kleiner als die der geladenen Fermionen sein sollten.
* Der Higgs-Sektor des SM ist, sowohl was die Anzahl der Higgs-Teilchen als auch die Einordnung in Darstellungen von $`G_{\text{SM}}`$ angeht, weitgehend willkürlich. Auch die Selbstkopplung sowie die Yukawa-Kopplungen und Massen der Higgs-Teilchen sind nicht festgelegt.
* Der QCD-$`\theta `$-Parameter ist extrem klein, $`\theta 10^9`$. Dies folgt direkt aus den experimentellen Grenzen für das Dipolmoment des Neutrons.
* Die durch Einsteins Allgemeine Relativitätstheorie beschriebene Gravitation läßt sich im Gegensatz zu den restlichen Wechselwirkungen nicht allein mit dem Prinzip der lokalen Eichinvarianz erklären. Alle bisher entwickelten Eichtheorien der Gravitation haben sich als nichtrenormierbar erwiesen.
Zahlreiche Modifikationen des minimalen SM sind konstruiert worden, um einzelne dieser Schwächen zu beheben. So kann das Problem der $`CP`$-Verletzung in der QCD durch ein Modell mit zwei Higgs-Doubletts und einer zusätzlichen chiralen $`U(1)`$-Symmetrie gelöst werden , während Modelle mit sogenannten horizontalen Symmetrien versuchen, die Massenhierarchie zwischen den Fermionfamilien zu erklären .
Im Gegensatz dazu gehen die Grand Unified-Theorien (GUTs), die im nächsten Kapitel diskutiert werden, über eine bloße Erweiterung des SM hinaus. Es wird vielmehr der Versuch gemacht, das SM in eine umfassendere Theorie einzubetten. GUTs können einige, wenn auch nicht alle, der oben erwähnten Schwächen beseitigen und zum besseren Verständnis des SM beitragen.
## Kapitel 2 Grand Unified-Theorien
### 2.1 Grundidee und allgemeine Eigenschaften
Das Hauptziel bei der Konstruktion von Grand Unified-Theorien (GUTs) besteht darin, die drei qualitativ und quantitativ verschiedenen Wechselwirkungen des SM in einer einzigen zu vereinheitlichen. Formal geschieht dies durch die Einbettung der SM-Symmetriegruppe $`G_{\text{SM}}`$ in eine einfache Lie-Gruppe $`GG_{\text{SM}}`$. Diese neue Theorie soll das SM als effektive Niederenergienäherung enthalten, was durch spontane Symmetriebrechung in einem oder auch mehreren Schritten erreicht werden kann.
Motiviert wird dieser Ansatz in erster Linie durch die Skalenabhängigkeit der Eichkopplungen im SM . Integriert man die Renormierungsgruppengleichungen (B.7-B.9) für die drei Größen $`\alpha _i=g_i^2/4\pi `$ von der Skala $`M_Z`$ zu höheren Energien, erhält man das in Abbildung 2.1 gezeigte Verhalten.
Die laufenden Kopplungen treffen sich näherungsweise im Bereich $`\mu 10^{14}`$-$`10^{15}\text{GeV}`$ und $`\alpha _i1/40`$. Daß sie sich jedoch nachweislich nicht in einem Punkt kreuzen, ist erst seit der Inbetriebnahme des Teilchenbeschleunigers LEP am CERN 1989 und der damit verbundenen Steigerung der Meßwertgenauigkeit für $`\alpha _i(M_Z)`$ bekannt .
Der erste Schritt bei der Konstruktion eines GUT-Modells besteht in der Wahl einer kompakten und einfachen Lie-Gruppe $`G`$. Man kann prinzipiell auch halbeinfache Lie-Gruppen, also direkte Produkte von einfachen Gruppen, verwenden, erhält dann aber für jeden Faktor eine separate Kopplungskonstante. Damit diese Kopplungen gleich sind und eine Vereinheitlichung zustande kommt, müssen zusätzliche Symmetrien eingeführt werden. Deshalb sind einfache Gruppen vorzuziehen, da man auf diese Weise eine Wechselwirkung mit genau einer Kopplung erhält; die SM-Wechselwirkungen sind dann verschiedene Aspekte dieser Kraft.
Die einfachen Lie-Gruppen beziehungsweise -algebren sind vollständig klassifiziert und ihre Eigenschaften und die ihrer Darstellungen bekannt ; Tabelle 2.1 gibt eine Übersicht.
Es gibt neben den vier unendlichen Reihen der klassischen Lie-Gruppen die fünf exzeptionellen Gruppen. Der Rang ist gleich der maximalen Anzahl miteinander kommutierender Generatoren, die Ordnung ist die Gesamtzahl der Generatoren und somit die Dimension der adjungierten Darstellung, und Dim(f) gibt die Dimension der Fundamentaldarstellung an.
Um eine geeignete Symmetriegruppe $`G`$ zu finden, sollte diese einer Reihe von Anforderungen genügen:
* Zunächst muß $`G`$ die SM-Symmetriegruppe $`G_{\text{SM}}`$ als Untergruppe enthalten. Da der Rang von $`SU(3)_CSU(2)_LU(1)_Y`$ gleich 4 ist, muß auch $`G`$ mindestens vom Rang 4 sein. Andererseits ist die Forderung nach der Minimalität des Ranges der Eichgruppe sicher sinnvoll, da mit dem Rang auch die Dimension der Darstellungen und somit der Umfang des Teilchenspektrums zunimmt.
* Die Gruppe $`G`$ sollte komplexe Darstellungen besitzen (eine Darstellung $``$ heißt komplex, wenn sie zu ihrer komplex konjugierten Darstellung $`\overline{}`$ nicht äquivalent ist; ansonsten wird sie als reell bezeichnet). Diese Forderung ist aus folgendem Grund plausibel:
Sei $`_{L,R}`$ die Darstellung, nach der sich die links- beziehungsweise rechtshändigen Teilchen und Antiteilchen einer Fermionfamilie transformieren. Dann gilt wegen $`\mathrm{\Psi }_R=C(\overline{\mathrm{\Psi }}_L^C)^T`$ die Beziehung $`_R=\overline{}_L`$. Im SM ist die $`_L`$ entsprechende Darstellung (1.5) komplex, was die Paritätsverletzung und die Nichtexistenz $`G_{\text{SM}}`$-invarianter Massenterme zur Folge hat. Um im Rahmen eines GUT-Modells eine reelle Darstellung $`_L`$ verwenden zu können, müssen in ihr neben den SM-Fermionen zusätzlich linkshändige Fermionen liegen, die sich wie das komplex konjugierte von (1.5) transformieren. Diese neuen Teilchen werden auch als Spiegelfermionen bezeichnet. Da sie in der Natur jedoch nicht beobachtet werden, müssen ihre Massen während des Symmetriebrechungsschrittes entstehen, der auf das SM führt, was aufgrund von Mischungen wiederum superschwere Massen auch für die SM-Fermionen zur Folge hat. Dieses Problem läßt sich auf einfache Weise durch Verwendung komplexer Darstellungen für die Fermionen vermeiden. Von den einfachen Lie-Gruppen enthalten jedoch nur
$$SU(n)(n3),SO(4n+2)(n2)\text{und}E_6$$
(2.1)
komplexe Darstellungen , was die Wahlmöglichkeiten für $`G`$ stark einschränkt.
* Damit die Renormierbarkeit der Theorie sichergestellt ist, muß sie anomaliefrei sein, also (1.7) erfüllen. Während die orthogonalen Gruppen (außer $`SO(6)SU(4)`$) und $`E_6`$ automatisch anomaliefrei sind, gilt dies bei den $`SU(n)`$-Gruppen nur für bestimmte Kombinationen irreduzibler Darstellungen . Im einfachsten Fall, der hier von Interesse ist, nämlich $`SU(5)`$, ist die (komplexe) Summe
$$\overline{\mathrm{𝟓}}\mathbf{\hspace{0.33em}10}\left[\mathrm{𝟏}\mathbf{\hspace{0.33em}16}_{SO(10)}\right]$$
(2.2)
anomaliefrei.
Berücksichtigt man bei der Wahl von $`G`$ diese Kriterien, so verbleiben folgende Möglichkeiten mit 4 $``$ Rang($`G`$) $``$ 6:
$$\text{Rang 4:}SU(5);\text{Rang 5:}SU(6),SO(10);\text{Rang 6:}SU(7),E_6$$
(2.3)
Auf $`SU(5)`$, $`SO(10)`$ und $`E_6`$ basierende GUTs werden in den nächsten Abschnitten behandelt; $`SU(6)`$\- und $`SU(7)`$-Theorien bieten gegenüber diesen Modellen keinerlei Vorteile, sind aber in einigen Punkten wesentlich unhandlicher und deswegen weitgehend unbeachtet geblieben.
Hat man auf diese Weise eine Wahl für die Eichgruppe $`G`$ getroffen, so sind noch folgende Schritte durchzuführen:
* Die Einbettung von $`G_{\text{SM}}`$ in $`G`$ und das Schema
$$G\stackrel{M_U}{}G_1\stackrel{M_I^{(1)}}{}\mathrm{}\stackrel{M_I^{(k)}}{}G_{\text{SM}}\stackrel{M_Z}{}SU(3)_CU(1)_{\text{em}}$$
(2.4)
für die Symmetriebrechung in das SM sind anzugeben. Der einfachste Fall, die direkte Symmetriebrechung $`G\stackrel{M_U}{}G_{\text{SM}}`$, kommt allerdings wegen der Tatsache, daß die laufenden SM-Kopplungen sich nicht in einem Punkt treffen (siehe Abbildung 2.1), nicht in Frage.
* Für jeden Symmetriebrechungsschritt in (2.4) müssen geeignete Higgs-Darstellungen festgelegt und zugehörige Potentiale konstruiert werden, welche ihn realisieren können. Dabei ist darauf zu achten, daß die Teilchen des SM nach wie vor erst im letzten Schritt bei $`M_Z`$ ihre Massen erhalten, während alle zusätzlichen Teilchen sehr viel schwerer sind und deshalb mindestens Massen der Größenordnung $`M_I^{\left(k\right)}`$ haben. Dieser Sachverhalt wird auch als „Survival Hypothesis“ bezeichnet.
* Die komplexe und anomaliefreie Darstellung, in der die Fermionen liegen sollen, ist auszuwählen. Hierbei muß die Zerlegung der Fermiondarstellung unter $`G_{\text{SM}}`$ die Zuordnung der korrekten SM-Quantenzahlen gestatten, das heißt die SM-Darstellung (1.5) enthalten.
Ist das geschehen, so ist die formale Konstruktion der GUT abgeschlossen. Die expliziten Eigenschaften und Vorhersagen des Modells, wie zum Beispiel Vereinheitlichungsmasse $`M_U`$ und -kopplung $`\alpha _U`$, Fermionmassenrelationen oder die Rate des Protonzerfalls und daraus eventuell resultierende Widersprüche zu experimentellen Resultaten sind dann im Detail zu untersuchen.
### 2.2 $`SU(5)`$
#### 2.2.1 Konstruktion und Teilcheninhalt
Die $`SU(5)`$-GUT stellt gewissermaßen den Prototyp dieser Theorien dar. Wie oben gezeigt wurde, ist sie die einzige für die Konstruktion einer GUT in Frage kommende Gruppe mit Rang($`G`$)=4. Da die SM-Eichgruppe $`G_{\text{SM}}`$ eine maximale Untergruppe der $`SU(5)`$ ist, kann nur die direkte Symmetriebrechung
$$SU(5)\stackrel{M_U}{}G_{\text{SM}}\stackrel{M_Z}{}SU(3)_CU(1)_{\text{em}}$$
(2.5)
erfolgen, was wegen des durch Abbildung 2.1 verdeutlichten SM-Kopplungsverhaltens für eine realistische Theorie nicht in Frage kommt. Dennoch sollen hier kurz die grundlegenden Eigenschaften der $`SU(5)`$-GUT zusammengefaßt werden, da sie sich zum großen Teil zumindest qualitativ auch auf andere GUT-Modelle übertragen lassen.
Betrachtet man die Verzweigungen der beiden $`SU(5)`$-Darstellungen $`\overline{\mathrm{𝟓}}`$ und $`\mathrm{𝟏𝟎}`$ unter $`G_{\text{SM}}`$
$$\overline{\mathrm{𝟓}}(\overline{\mathrm{𝟑}}\mathbf{,1})_{{\scriptscriptstyle \frac{1}{3}}}(\mathbf{1,2})_{{\scriptscriptstyle \frac{1}{2}}},\mathrm{𝟏𝟎}(\mathbf{3,2})_{{\scriptscriptstyle \frac{1}{6}}}(\overline{\mathrm{𝟑}}\mathbf{,1})_{{\scriptscriptstyle \frac{2}{3}}}(\mathbf{1,1})_1,$$
(2.6)
stellt man fest, daß die reduzible Darstellung $`(\overline{\mathrm{𝟓}}\mathrm{𝟏𝟎})`$ gerade alle SM-Fermionen einer Familie aufnehmen kann (siehe Tabelle 1.1). Ferner ist sie komplex und anomaliefrei, erfüllt also alle notwendigen Bedingungen. Die gemeinsame Einbettung von Quarks und Leptonen in dieselbe irreduzible Darstellung ist eine Eigenschaft aller GUTs und führt, da die Eichtransformationen diese Teilchen miteinander mischen können, zur Verletzung der Baryon- und Leptonzahlerhaltung.
An dieser Stelle wird ein wesentlicher Erfolg von GUTs deutlich, nämlich die Erklärung für die Quantisierung der elektrischen Ladung. Sie gehört zu den allgemeinen Eigenschaften von Eichtheorien mit einfacher Symmetriegruppe, da die Eigenwerte der diagonalen Generatoren solcher Gruppen im Gegensatz zu den Eigenwerten der abelschen $`U(1)`$ stets diskret sind. Der Ladungsoperator $`Q`$ muß in der Cartan-Unteralgebra von $`G`$ liegen, also eine Linearkombination der diagonalen Generatoren sein. Da die Generatoren spurfrei sind, gilt dies auch für $`Q`$. Daraus folgt wiederum, daß die Summe der Ladungen der in den Darstellungen $`\overline{\mathrm{𝟓}}`$ und $`\mathrm{𝟏𝟎}`$ liegenden Fermionen jeweils 0 sein muß: $`Q(e^{})=3Q(d)`$ und $`2Q(u)+Q(e^{})+Q(d)=0`$. Quarks haben deshalb drittelzahlige Ladungen, weil sie in drei Farben auftreten, während Leptonen farbneutral sind.
Betrachtet man $`G_{\text{SM}}`$ als Untergruppe von $`SU(5)`$, sind die Hyperladungswerte in Tabelle 1.1 mit einem Faktor $`\sqrt{5/3}`$ umzunormieren. Wegen (1.22), $`g^{}=\sqrt{3/5}g_1`$ und $`g_1(M_U)=g_2(M_U)`$ gilt in GUTs mit einfacher Symmetriegruppe stets $`\mathrm{sin}^2\theta _W(M_U)=3/8`$; der Weinbergwinkel ist kein freier Parameter mehr. Der experimentelle Wert $`\mathrm{sin}^2\theta _W(M_Z)0.23`$ muß durch Integration der entsprechenden Renormierungsgruppengleichungen reproduziert werden.
Die adjungierte Darstellung, nach der sich in Eichtheorien die Eichbosonen transformieren, verzweigt sich bezüglich $`G_{\text{SM}}`$ gemäß
$$\mathrm{𝟐𝟒}(\mathbf{8,1})_0(\mathbf{1,3})_0(\mathbf{1,1})_0(\mathbf{3,2})_{{\scriptscriptstyle \frac{5}{6}}}(\overline{\mathrm{𝟑}}\mathbf{,2})_{{\scriptscriptstyle \frac{5}{6}}}$$
(2.7)
Die ersten drei Summanden entsprechen den SM-Eichbosonen, während die Darstellungen $`(\mathbf{3,2})_{{\scriptscriptstyle \frac{5}{6}}}`$ und $`(\overline{\mathrm{𝟑}}\mathbf{,2})_{{\scriptscriptstyle \frac{5}{6}}}`$ neue Bosonen enthalten, die mit $`X`$ und $`Y`$ beziehungsweise $`\overline{X}`$ und $`\overline{Y}`$ bezeichnet werden. Da diese experimentell nicht beobachtet werden, müssen sie Massen der Größenordnung $`M_U`$ haben.
#### 2.2.2 Protonzerfall
Die $`X`$\- und $`Y`$-Bosonen, welche die elektrischen Ladungen +4/3 und +1/3 besitzen, tragen sowohl eine Farbladung als auch schwachen Isospin und können deshalb Quarks und Leptonen ineinander überführen. Während Baryon- und Leptonzahl $`B,L`$ globale $`U(1)`$-Symmetrien des (perturbativen) SM sind, werden sie in der $`SU(5)`$-Theorie wie auch in allen anderen GUTs verletzt; allerdings bleibt $`BL`$ erhalten. Die direkte Konsequenz der Baryonzahlverletzung ist die Instabilität von Proton und (gebundenem) Neutron. Der Nukleonenzerfall ist eine wichtige und experimentell überprüfbare Vorhersage von GUTs, die über das SM klar hinausgeht .
Der baryonzahlverletzende Teil der Lagrangedichte für die erste Fermionfamilie (unter Vernachlässigung von Mischungen) ergibt sich zu :
$`_{\mathrm{\Delta }B0}`$ $`=`$ $`{\displaystyle \frac{g_U}{\sqrt{2}}}\overline{X}_\mu ^\alpha \left(\epsilon _{\alpha \beta \gamma }\overline{u}_L^{C\gamma }\gamma ^\mu u_L^\beta +\overline{d}_{L\alpha }\gamma ^\mu e_L^++\overline{d}_{R\alpha }\gamma ^\mu e_R^+\right)`$ (2.8)
$`+`$ $`{\displaystyle \frac{g_U}{\sqrt{2}}}\overline{Y}_\mu ^\alpha \left(\epsilon _{\alpha \beta \gamma }\overline{u}_L^{C\gamma }\gamma ^\mu d_L^\beta \overline{d}_{R\alpha }\gamma ^\mu \nu _{eR}^C\overline{u}_{L\alpha }\gamma ^\mu e_L^+\right)+\text{h.c.}`$
$`\alpha `$, $`\beta `$ und $`\gamma `$ sind $`SU(3)_C`$-Farbindizes. Das entspricht den in Abbildung 2.2 aufgeführten Wechselwirkungsvertizes.
Zwei typische Beispiele für Feynmandiagramme von Protonzerfallsprozessen zeigt Abbildung 2.3:
Der zugrundeliegende Prozeß lautet stets Quark+Quark $``$ Antiquark+Antilepton, was zu Zerfallskanälen Nukleon $``$ Meson+Antilepton führt. Alle anderen Elementarprozesse sind aufgrund von Phasenraumeffekten unterdrückt. Unter der Annahme kleiner Mischungen ist $`p\pi ^0e^+`$ der dominante Kanal; eine qualitative Abschätzung der Lebensdauer gemäß
$$\mathrm{\Gamma }_p=\tau _p^1\alpha _U^2\frac{m_p^5}{M_{X,Y}^4}$$
(2.9)
liefert mit $`\alpha _U1/40`$ und $`M_{X,Y}M_U510^{14}`$ GeV für $`\tau _p`$ den Wert $`310^{30}`$ Jahre. Die neuesten experimentellen Grenzen liegen bei $`\tau _p2.910^{33}`$ Jahre . Die zu kurze Lebensdauer des Protons war die erste Vorhersage der $`SU(5)`$-GUT, die eindeutig im Widerspruch zu experimentellen Resultaten stand.
#### 2.2.3 Symmetriebrechung und Fermionmassen
Um den ersten Symmetriebrechungsschritt in (2.5) zu verwirklichen, benötigt man eine Higgs-Darstellung $`\mathrm{\Phi }`$, die ein SM-Singulett enthält; dieses bildet dann den Vakuumerwartungswert (VEW) aus. Der zweite Brechungsschritt bei $`M_Z`$ erfordert eine Darstellung $`H`$, in der ein SM-Doublett $`(\mathbf{1,2})_{{\scriptscriptstyle \frac{1}{2}}}`$ liegt. Die einfachste Wahl hierfür lautet $`\mathrm{\Phi }=\mathrm{𝟐𝟒}`$ und $`H=\mathrm{𝟓}`$ (siehe (2.6) und (2.7)). Der VEW $`\mathrm{\Phi }`$ gibt allen Teilchen, die nicht zum SM-Spektrum gehören, Massen der Größenordnung $`M_U`$, während $`H`$ die SM-Massen erzeugt.
Betrachtet man das Transformationsverhalten von Fermionmassentermen in der $`SU(5)`$-Theorie, so ergeben sich die drei Tensorprodukte
$`\overline{\mathrm{𝟓}}\overline{\mathrm{𝟓}}`$ $`=`$ $`\overline{\mathrm{𝟏𝟎}}\overline{\mathrm{𝟏𝟓}}`$ (2.10)
$`\overline{\mathrm{𝟓}}\mathbf{\hspace{0.17em}10}`$ $`=`$ $`\mathrm{𝟓}\mathbf{\hspace{0.17em}45}`$ (2.11)
$`\mathrm{𝟏𝟎}\mathbf{\hspace{0.17em}10}`$ $`=`$ $`\overline{\mathrm{𝟓}}\overline{\mathrm{𝟒𝟓}}\overline{\mathrm{𝟓𝟎}}`$ (2.12)
Daran erkennt man zunächst, daß keine Fermionmassen $`\mathrm{\Phi }`$ entstehen können. Die Kopplung von $`H`$ an den Term $`\overline{\mathrm{𝟓}}\mathrm{𝟏𝟎}`$ erzeugt Dirac-Massen für die $`d`$-Quarks und die geladenen Leptonen, jene an den Term $`\mathrm{𝟏𝟎}\mathrm{𝟏𝟎}`$ Massen für die $`u`$-Quarks. Majorana-Massen für die Neutrinos entstehen nicht, da $`H`$ nicht an $`\overline{\mathrm{𝟓}}\overline{\mathrm{𝟓}}`$ koppelt.
Bei $`M_U`$ gilt die Beziehung $`𝐌_d=𝐌_e`$, während $`𝐌_u`$ von den anderen Massenmatrizen unabhängig ist. Vorhersagen von Relationen zwischen den fermionischen Massenmatrizen gehören zu den typischen Eigenschaften von GUTs. In diesem Fall ist die Vorhersage $`𝐌_d(M_U)=𝐌_e(M_U)`$ nicht realistisch, da sich nach Berücksichtigung der Renormierungseffekte die Beziehung $`𝐌_d(M_Z)3𝐌_e(M_Z)`$ ergibt, was im Widerspruch zu den Fermionmassen aus Tabelle 1.6 steht. Durch Einführung einer zusätzlichen 45-Higgs-Darstellung, die wegen (2.102.12) ebenfalls Fermionmassen erzeugen kann, ist es möglich, die Relation zwischen $`𝐌_d`$ und $`𝐌_e`$ so zu modifizieren, daß sie zumindest qualitativ korrekte Ergebnisse liefert .
#### 2.2.4 Das Hierarchieproblem
Die $`SU(5)`$-GUT enthält zwei verschiedene Symmetriebrechungsskalen $`M_U`$ und $`M_Z`$. Während die SM-Teilchen nach wie vor Massen der Größenordnung $`M_Z`$ haben sollen, müssen die zusätzlich eingeführten Teilchen Massen $`M_U`$ besitzen. Das Verhältnis dieser Skalen ist extrem klein, $`M_Z/M_U10^{12}`$. Diese Hierarchie überträgt sich zwangsläufig auf die Parameter im Higgs-Potential $`V`$, welches die VEW und damit die Symmetriebrechung verursacht . Das vollständige Higgs-Potential beinhaltet neben
$`V(\mathrm{\Phi })`$ $`=`$ $`m_1^2\text{Tr}\mathrm{\Phi }^2+\lambda _1\left(\text{Tr}\mathrm{\Phi }^2\right)^2+\lambda _2\text{Tr}\mathrm{\Phi }^4\text{und}`$ (2.13)
$`V(H)`$ $`=`$ $`m_2^2H^{}H+\lambda _3(H^{}H)^2`$ (2.14)
auch Mischterme der Form
$$V(\mathrm{\Phi },H)=\lambda _4\text{Tr}\mathrm{\Phi }^2(H^{}H)+\lambda _5\left(H^{}\mathrm{\Phi }^2H\right)$$
(2.15)
Der erste Brechungsschritt erfolgt, wenn $`\mathrm{\Phi }`$ den VEW
$$\mathrm{\Phi }=\upsilon _1\text{diag}(\mathrm{2,2,2},3,3)\text{mit}\upsilon _1^2=m_1^2/(60\lambda _1+14\lambda _2)M_U^2$$
(2.16)
ausbildet. Die $`X,Y`$-Bosonen und die physikalischen Higgs-Teilchen in $`\mathrm{\Phi }`$ erhalten Massen $`M_U`$. Der zweite Schritt verläuft analog zum SM. Berechnet man jedoch die Massen der Teilchen in $`H`$, so ergibt sich aufgrund der Potentialterme in (2.15):
$$(\mathbf{3,1})_H:m_T^2=m_2^2+(30\lambda _4+4\lambda _5)\upsilon _1^2,(\mathbf{1,2})_H:m_D^2=m_2^2+(30\lambda _4+9\lambda _5)\upsilon _1^2$$
(2.17)
Sowohl $`m_T`$ als auch $`m_D`$ enthalten Terme $`\upsilon _1`$, was im Falle des Tripletts $`(\mathbf{3,1})_H`$ auch notwendig ist, da es ebenso wie die $`X,Y`$-Bosonen Protonzerfälle vermitteln kann; das SM-Doublett $`(\mathbf{1,2})_H`$ dagegen muß Massen $`M_Z`$ haben. Dies läßt sich wegen $`|m_2/\upsilon _1|10^{12}`$ nur dann erreichen, wenn $`(30\lambda _4+9\lambda _5)10^{24}`$ ist.
Die nicht sehr natürlich anmutende Feineinstellung der Parameter $`\lambda _{\mathrm{4,5}}`$ mit einer Genauigkeit $`10^{24}`$ muß schließlich für jede Ordnung der Störungstheorie wiederholt werden, da sonst Strahlungskorrekturen wie zum Beispiel die in Abbildung 2.4 gezeigte Beiträge höherer Ordnung zu $`m_D`$ liefern, die wieder $`\upsilon _1`$ sind .
Dieses Hierarchieproblem taucht grundsätzlich in allen Theorien auf, die zwei oder mehr Massenskalen stark unterschiedlicher Größenordnung aufweisen und läßt sich zumindest im Rahmen nichtsupersymmetrischer GUTs nicht befriedigend lösen. Deswegen geht man im allgemeinen von der Gültigkeit der sogenannten „Extended Survival Hypothesis“ aus . Sie besagt, daß Higgs-Teilchen die größtmögliche Masse erhalten, die mit dem Symmetriebrechungsschema des Modells verträglich ist. Da in der $`SU(5)`$-GUT das Triplett $`(\mathbf{3,1})_H`$ für den Brechungsschritt bei $`M_Z`$ ohne Bedeutung ist, hat es eine Masse $`M_U`$.
#### 2.2.5 Zusammenfassung
Obwohl das $`SU(5)`$-Modell aufgrund der experimentellen Fakten mittlerweile ausgeschlossen ist, besitzt es im Vergleich zum SM eine Reihe von attraktiven Eigenschaften:
* Die SM-Wechselwirkungen werden vereinheitlicht; es gibt nur noch eine Kopplungskonstante $`g_U`$ und eine Vorhersage für $`\mathrm{sin}^2\theta _W(M_U)`$.
* Die anomaliefreie Darstellung $`(\overline{\mathrm{𝟓}}\mathrm{𝟏𝟎})`$ kann die Fermionen einer Familie aufnehmen. Obwohl nach wie vor reduzibel, hat sie eine einfachere Struktur als (1.5).
* Die Einbettung der $`U(1)_Y`$ in eine einfache Lie-Gruppe führt automatisch zur korrekten Quantisierung der elektrischen Ladung.
* Die Massenmatrizen der Fermionen sind nicht mehr unabhängig voneinander.
Dem stehen folgende Schwierigkeiten gegenüber:
* Die fermionische Darstellung ist reduzibel.
* Der Higgs-Sektor der Theorie ist wesentlich komplexer als im SM.
* Die Existenz zweier stark unterschiedlicher Massenskalen führt zum Hierarchieproblem.
* Die Probleme im Neutrinosektor sind weiterhin ungelöst.
* Auch in der $`SU(5)`$-GUT ist das Wechselwirkungsverhalten rechts- und linkshändiger Fermionen verschieden.
* Es gibt keine Erklärung für die Existenz dreier Fermionfamilien und die Massenhierarchie zwischen ihnen.
Einige dieser Schwächen können im Rahmen von GUTs mit größerer Eichgruppe beseitigt werden.
### 2.3 $`SO(10)`$
#### 2.3.1 Spinordarstellungen der $`SO(n)`$-Gruppen
Die $`SO(10)`$ ist die kleinste orthogonale Gruppe, die $`G_{\text{SM}}`$ als Untergruppe enthält und komplexe Darstellungen besitzt. Da $`SU(5)U(1)`$ eine maximale Untergruppe der $`SO(10)`$ ist, sind alle vorteilhaften Eigenschaften der $`SU(5)`$-GUT auch in $`SO(10)`$-Theorien vorhanden.
Die orthogonalen Lie-Gruppen unterscheiden sich in einigen wichtigen Punkten von den unitären Gruppen . Alle irreduziblen Darstellungen der $`SU(n)`$ können durch Reduktion von Tensorprodukten der komplexen Fundamentaldarstellung $`𝐧`$ und der zu ihr komplex konjugierten Darstellung $`\overline{𝐧}`$ erhalten werden. Die orthogonalen Gruppen $`SO(n)`$ dagegen besitzen neben den Tensordarstellungen, welche man durch Produktbildung aus der reellen Fundamentaldarstellung $`𝐧`$ konstruieren kann, die sogenannten Spinordarstellungen, für die das nicht gilt. Zu den Spinordarstellungen der $`SO(n)`$ gelangt man über die komplexe Clifford-Algebra
$$C_n:\{\mathrm{\Gamma }_j,\mathrm{\Gamma }_k\}=\mathrm{\hspace{0.33em}2}\delta _{jk}\mathbf{\hspace{0.17em}1}(j,k=1,\mathrm{},n)$$
(2.18)
Aus den Generatoren von $`C_n`$ kann man eine Darstellung für die Generatoren der $`SO(n)`$-Lie-Algebra konstruieren:
$$\mathrm{\Sigma }_{jk}=\frac{i}{4}[\mathrm{\Gamma }_j,\mathrm{\Gamma }_k]$$
(2.19)
Die so definierten $`\mathrm{\Sigma }_{jk}`$ erfüllen die Vertauschungsrelation
$$[_{ij},_{kl}]=i\left(\delta _{ik}_{jl}\delta _{il}_{jk}\delta _{jk}_{il}+\delta _{jl}_{ik}\right),$$
(2.20)
welche die $`SO(n)`$-Algebra definiert ($`_{ij}=_{ji}\text{mit}i,j=1,\mathrm{},n`$). Die Generatoren $`_{ij}`$ in der Fundamentaldarstellung sind durch
$$\left(_{ij}\right)_{ab}=i(\delta _{ia}\delta _{jb}\delta _{ja}\delta _{ib})$$
(2.21)
gegeben. Über die Matrixdarstellungen der komplexen Clifford-Algebra (2.18) und somit der $`\mathrm{\Sigma }_{jk}`$ kommt man zu den Spinoren, da diese die Elemente des entsprechenden Darstellungsraumes sind. Die Eigenschaften der irreduziblen Spinordarstellungen der $`SO(n)`$ hängen vom Wert von $`n`$ ab:
* Ist $`n`$ ungerade, so existiert eine reelle Spinordarstellung der Dimension $`2^{(n1)/2}`$.
* Ist $`n`$ gerade, so existieren zwei inäquivalente Spinordarstellungen der Dimension $`2^{(n/21)}`$. Wenn $`n/2`$ gerade ist, so sind die beiden Darstellungen reell, während sie für ungerades $`n/2`$ komplex sind, wobei dann die eine das komplex Konjugierte der anderen ist.
Aus diesem Grunde verfügen nur die $`SO(4n+2)`$-Gruppen über komplexe Darstellungen; im Falle der $`SO(10)`$ sind dies die $`\mathrm{𝟏𝟔}`$ und die $`\overline{\mathrm{𝟏𝟔}}`$.
#### 2.3.2 Teilcheninhalt und Symmetriebrechung
Die Hauptmotivation für die Konstruktion von $`SO(10)`$-GUTs wird deutlich, wenn man die Verzweigung der 16 bezüglich der Untergruppe $`SU(5)`$ betrachtet:
$$\mathrm{𝟏𝟔}\mathrm{𝟏𝟎}\overline{\mathrm{𝟓}}\mathrm{𝟏}$$
(2.22)
In einer 16 können demnach alle Fermionen einer Familie und ein zusätzliches $`SU(5)`$-Singulett untergebracht werden. Dieses Singulett besitzt exakt die Eigenschaften und Quantenzahlen des rechtshändigen Neutrinos (beziehungsweise linkshändigen Antineutrinos), da es zur $`G_{\text{SM}}`$-Darstellung $`(\mathbf{1,1})_0`$ identisch ist.
Hier wird ein großer Vorteil von $`SO(10)`$-Modellen gegenüber der $`SU(5)`$-GUT deutlich, denn alle Fermionen einer Generation inklusive eines linkshändigen Antineutrinos liegen in einer irreduziblen und komplexen Darstellung der Eichgruppe. Ferner erklärt die Anomaliefreiheit der orthogonalen Gruppen auf natürliche Weise das Verschwinden der Anomalie der $`SU(5)`$-Darstellung $`\mathrm{𝟏𝟎}\overline{\mathrm{𝟓}}`$.
$`SO(10)`$-Modelle sind überdies zumindest im Eichboson- und Fermionsektor manifest $`C`$\- und $`P`$-invariant. Da in der 16 die linkshändigen Teilchen und Antiteilchen liegen, wird sie durch Anwendung von $`C`$ auf sich selbst abgebildet. Unter $`P`$ wird die 16 auf die $`\overline{\mathrm{𝟏𝟔}}`$ abgebildet, in welcher die rechtshändigen Teilchen und Antiteilchen untergebracht sind. $`C`$\- und $`P`$-Verletzung kann entweder explizit im Higgs-Sektor oder durch die spontane Symmetriebrechung realisiert werden.
Da der Rang der $`SO(10)`$ um eins größer als der von $`G_{\text{SM}}`$ ist, existiert eine Vielzahl von möglichen Symmetriebrechungsschemata. Die direkte Brechung in das SM ist dabei ebenso wie die Brechung nach $`SU(5)[U(1)]`$ zwar möglich, aber wegen der in Abbildung 2.1 gezeigten SM-Kopplungsentwicklung phänomenologisch ausgeschlossen. Man benötigt demnach mindestens einen Zwischenschritt, das heißt
$$SO(10)\stackrel{M_U}{}G_I\stackrel{M_I}{}G_{\text{SM}}\stackrel{M_Z}{}SU(3)_CU(1)_{\text{em}}$$
(2.23)
Es sind Modelle mit bis zu vier aufeinander folgenden intermediären Symmetriegruppen konstruiert und ausführlich untersucht worden , jedoch besitzen sie gegenüber Modellen mit nur einer solchen keine besonderen Vorzüge. Tabelle 2.2 gibt die verschiedenen Möglichkeiten für den Brechungsschritt $`SO(10)G_I`$ bei $`M_U`$ und die dafür zu verwendende Higss-Darstellung mit einem $`G_I`$-Singulett an.
Die $`D`$-Parität ist eine diskrete Symmetrie, welche $`SU(2)_L`$ und $`SU(2)_R`$ vertauscht; sie erfordert ein rechts-links-symmetrisches Teilchenspektrum. Wenn die Theorie also Teilchen in einer $`SU(2)_LSU(2)_R`$-Darstellung $`(m,n)`$ enthält, muß auch eine Darstellung $`(n,m)`$ vorhanden sein. Dies hat zur Folge, daß die beiden Eichkopplungen $`g_{2L}(\mu )`$ und $`g_{2R}(\mu )`$ überall zwischen $`M_U`$ und $`M_I`$ gleich groß sind und die Parität erhalten ist.
Welche Symmetriegruppe $`G_I`$ durch die Brechung mittels einer bestimmten Darstellung $`\mathrm{\Phi }`$ realisiert wird, hängt von den jeweiligen Werten der Parameter im Higgs-Potential $`V(\mathrm{\Phi })`$ ab . Die Symmetriebrechung bei $`M_I`$ in das SM erfolgt in allen Fällen über einen VEW des $`G_{\text{SM}}`$-Singuletts einer $`SO(10)`$-Darstellung 126.
In dieser Arbeit soll ein $`SO(10)`$-Modell mit intermediärer $`SU(4)_CSU(2)_LSU(2)_RG_{\text{PS}}`$-Symmetrie untersucht werden. Diese Gruppe gehört zu den maximalen Untergruppen der $`SO(10)`$ und ist von Pati und Salam für eine Erweiterung des SM durch eine rechts-links-symmetrische Theorie vorgeschlagen worden . Allerdings hat sich gezeigt, daß Modelle mit $`G_{\text{PS}}D`$-Symmetrie für $`M_U`$ den Wert $`10^{15}`$ GeV liefern, welcher auf eine zu kurze Protonlebensdauer führt . $`G_{\text{PS}}`$-Modelle ohne $`D`$-Parität besitzen diesen Nachteil nicht. Der $`SU(2)_L`$-Faktor ist identisch mit dem in $`G_{\text{SM}}`$, die $`SU(2)_R`$ ist das rechtshändige Gegenstück dazu, und die $`SU(4)_C`$ schließlich ist eine erweiterte Farbgruppe mit der $`BL`$-Quantenzahl als „vierter Farbe“. In $`SO(10)`$-basierten Modellen ist $`BL`$ bei Skalen $`\mu M_I`$ eine lokale $`U(1)`$-Symmetrie, also Teil der Eichgruppe, und keine globale Symmetrie wie in der $`SU(5)`$-GUT. Das hat zur Folge, daß auch Bosonen nichtverschwindendes $`BL`$ besitzen können.
Die Fermiondarstellung 16 verzweigt sich bezüglich $`G_{\text{PS}}`$ und $`SU(3)_CSU(2)_L`$ gemäß
$`\mathrm{𝟏𝟔}`$ $``$ $`(\mathbf{4,2,1})(\overline{\mathrm{𝟒}}\mathbf{,1,2})`$
$``$ $`\left(\begin{array}{cccc}u_1& u_2& u_3& \nu _e\\ d_1& d_2& d_3& e^{}\end{array}\right)_L\left(\begin{array}{cccc}d_1^C& d_2^C& d_3^C& e^+\\ u_1^C& u_2^C& u_3^C& \nu _e^C\end{array}\right)_L`$
$``$ $`(\mathbf{3,2})(\mathbf{1,2})2(\overline{\mathrm{𝟑}}\mathbf{,1})2(\mathbf{1,1})`$ (2.25)
Der Operator der elektrischen Ladung ist durch $`Q=T_{3L}+T_{3R}+(BL)/2`$ gegeben, die SM-Hyperladung durch $`Y=T_{3R}+(BL)/2`$. Die adjungierte Darstellung 45, in der die Eichbosonen liegen, verzweigt sich bezüglich $`G_{\text{PS}}`$ wie
$$\mathrm{𝟒𝟓}(\mathbf{15,1,1})(\mathbf{1,3,1})(\mathbf{1,1,3})(\mathbf{2,2,6})$$
(2.26)
Die ersten drei Faktoren entsprechen den $`G_{\text{PS}}`$-Eichbosonen. Die (2,2,6) enthält neben den aus der $`SU(5)`$-GUT bekannten $`X`$\- und $`Y`$-Bosonen die ebenfalls baryon- und leptonzahlverletzende Prozesse vermittelnden $`X^{}`$\- und $`Y^{}`$-Bosonen (und deren Antiteilchen); sie haben die Ladungen $`+2/3`$ beziehungsweise $`1/3`$ und Massen der Größenordnung $`M_U`$. Abbildung 2.5 zeigt die zugehörigen Wechselwirkungsvertizes. In der (15,1,1) liegen die acht Gluonen der QCD, ein an $`BL`$ koppelndes neutrales Boson und die $`SU(3)_C`$-Tripletts $`X_3`$ und $`\overline{X}_3`$ mit $`Q(X_3)=+2/3`$.
Der vollständige baryonzahlverletzende Teil der $`SO(10)`$-Lagrangedichte für die erste Fermionfamilie und ohne Mischungen lautet :
$`_{\mathrm{\Delta }B0}`$ $`=`$ $`{\displaystyle \frac{g_U}{\sqrt{2}}}\overline{X}_\mu ^\alpha \left(\epsilon _{\alpha \beta \gamma }\overline{u}_L^{C\gamma }\gamma ^\mu u_L^\beta +\overline{d}_{L\alpha }\gamma ^\mu e_L^++\overline{d}_{R\alpha }\gamma ^\mu e_R^+\right)`$
$`+`$ $`{\displaystyle \frac{g_U}{\sqrt{2}}}\overline{Y}_\mu ^\alpha \left(\epsilon _{\alpha \beta \gamma }\overline{u}_L^{C\gamma }\gamma ^\mu d_L^\beta \overline{d}_{R\alpha }\gamma ^\mu \nu _{eR}^C\overline{u}_{L\alpha }\gamma ^\mu e_L^+\right)`$
$`+`$ $`{\displaystyle \frac{g_U}{\sqrt{2}}}X_\mu ^{}_{}{}^{}\alpha \left(\epsilon _{\alpha \beta \gamma }\overline{d}_L^{C\gamma }\gamma ^\mu d_L^\beta \overline{u}_{L\alpha }\gamma ^\mu \nu _{eL}^C\overline{u}_{R\alpha }\gamma ^\mu \nu _{eR}^C\right)`$
$`+`$ $`{\displaystyle \frac{g_U}{\sqrt{2}}}Y_\mu ^{}_{}{}^{}\alpha \left(\epsilon _{\alpha \beta \gamma }\overline{d}_L^{C\gamma }\gamma ^\mu u_L^\beta \overline{d}_{L\alpha }\gamma ^\mu \nu _{eL}^C\overline{u}_{R\alpha }\gamma ^\mu e_R^+\right)`$
$`+`$ $`{\displaystyle \frac{g_U}{\sqrt{2}}}X_{3\mu }^\alpha \left(\overline{d}_{L\alpha }\gamma ^\mu e_L^{}+\overline{d}_{R\alpha }\gamma ^\mu e_R^{}+\overline{u}_{L\alpha }\gamma ^\mu \nu _{eL}+\overline{u}_{R\alpha }\gamma ^\mu \nu _{eR}\right)`$
$`+`$ h.c. (2.27)
Durch direkten $`X_3`$-Austausch können keine Nukleonenzerfälle vermittelt werden, was auch notwendig ist, da deren Massen der Größenordnung $`M_I`$ zu sehr großen Zerfallsraten führen würden. Durch $`X_3`$-$`X^{}`$-Mischung entstehen in höheren Ordnungen $`BL`$-verletzende Nukleonzerfälle, die aber gegenüber den Prozessen führender Ordnung durch zusätzliche Faktoren $`1/M_U`$ unterdrückt sind .
#### 2.3.3 Fermionmassen
Dirac-Massenterme für die geladenen Fermionen und die Neutrinos haben die Struktur $`m(\mathrm{\Psi }_L^C)^TC\mathrm{\Psi }_L`$ und besitzen demnach das $`SO(10)`$-Transformationsverhalten
$$\mathrm{𝟏𝟔}\mathrm{𝟏𝟔}=\mathbf{\hspace{0.33em}10}\mathrm{𝟏𝟐𝟎}\mathrm{𝟏𝟐𝟔}$$
(2.28)
Die Fermionmassen und die Symmetriebrechung $`G_{\text{SM}}\stackrel{M_Z}{}SU(3)_CU(1)_{\text{em}}`$ kommen also zustande, wenn die farblosen und elektrisch neutralen Komponenten von Higgs-Teilchen in den Darstellungen $`\overline{\mathrm{𝟏𝟎}}`$, $`\overline{\mathrm{𝟏𝟐𝟎}}`$ oder $`\overline{\mathrm{𝟏𝟐𝟔}}`$ einen VEW der Größenordnung $`M_Z`$ ausbilden. Diese Darstellungen verzweigen sich folgendermaßen unter $`G_{\text{PS}}`$ (siehe Anhang A.2):
$`\mathrm{𝟏𝟎}`$ $``$ $`(\mathbf{1,2,2})(\mathbf{6,1,1})`$ (2.29)
$`\mathrm{𝟏𝟐𝟎}`$ $``$ $`(\mathbf{1,2,2})(\mathbf{15,2,2})(\mathbf{6,3,1})(\mathbf{6,1,3})(\mathbf{10,1,1})(\overline{\mathrm{𝟏𝟎}}\mathbf{,1,1})`$ (2.30)
$`\mathrm{𝟏𝟐𝟔}`$ $``$ $`(\mathbf{15,2,2})(\mathbf{10,3,1})(\overline{\mathrm{𝟏𝟎}}\mathbf{,1,3})(\mathbf{6,1,1})`$ (2.31)
Da die linkshändigen Teilchen in (4,2,1) und die linkshändigen Antiteilchen in $`(\overline{\mathrm{𝟒}}\mathbf{,1,2})`$ liegen, wird die Erzeugung von Dirac-Massen wegen
$$(\mathbf{4,2,1})(\overline{\mathrm{𝟒}}\mathbf{,1,2})=(\mathbf{15,2,2})(\mathbf{1,2,2})$$
(2.32)
durch die $`SU(2)_L`$-Doubletts in den $`G_{\text{PS}}`$-Darstellungen (1,2,2) der 10 und 120 sowie denen in den (15,2,2) der 120 und 126 vermittelt. Ferner können aufgrund von
$`(\mathbf{4,2,1})(\mathbf{4,2,1})(\overline{\mathrm{𝟏𝟎}}\mathbf{,3,1})`$ $`=`$ $`(\mathbf{1,1,1})\mathrm{}\mathrm{}`$ (2.33)
$`(\overline{\mathrm{𝟒}}\mathbf{,1,2})(\overline{\mathrm{𝟒}}\mathbf{,1,2})(\mathbf{10,1,3})`$ $`=`$ $`(\mathbf{1,1,1})\mathrm{}\mathrm{}`$ (2.34)
über den Term $`\mathrm{𝟏𝟔}\mathrm{𝟏𝟔}\overline{\mathrm{𝟏𝟐𝟔}}`$ auch Majorana-Massen für die links- und rechtshändigen Neutrinos auftreten. Da der VEW des SM-Singuletts in (10,1,3) den Symmetriebrechungsschritt $`G_{\text{PS}}\stackrel{M_I}{}G_{\text{SM}}`$ realisiert, sind die Majorana-Massen der rechtshändigen Neutrinos zwangsläufig von der Größenordnung $`M_I`$.
Die verschiedenen Yukawa-Kopplungen und VEW der 10, 120 und 126 sind in Tabelle 2.3 zusammengefaßt (die Indizes numerieren mehrere Darstellungen einer Art, wenn vorhanden).
Damit die beiden VEW, die zu jeder $`G_{\text{PS}}`$-Darstellung gehören, verschieden sein können, muß diese komplex sein. Das läßt sich durch Kombination zweier reeller Darstellungen $`_1,_2`$ zu $`=_1+i_2`$ erreichen.
Mit den in Tabelle 2.3 festgelegten Bezeichnungen für die Kopplungen und VEW (unter Vernachlässigung der oberen Indizes) ergeben sich für die Fermionmassenmatrizen in einer $`SO(10)`$-GUT folgende Identitäten :
$`𝐌_d`$ $`=`$ $`\upsilon _d𝐘_{\mathrm{𝟏𝟎}}+\omega _d𝐘_{\mathrm{𝟏𝟐𝟔}}+(\stackrel{~}{\upsilon }_d+\stackrel{~}{\omega }_d)𝐘_{\mathrm{𝟏𝟐𝟎}}`$ (2.35)
$`𝐌_e`$ $`=`$ $`\upsilon _d𝐘_{\mathrm{𝟏𝟎}}3\omega _d𝐘_{\mathrm{𝟏𝟐𝟔}}+(\stackrel{~}{\upsilon }_d3\stackrel{~}{\omega }_d)𝐘_{\mathrm{𝟏𝟐𝟎}}`$ (2.36)
$`𝐌_u`$ $`=`$ $`\upsilon _u𝐘_{\mathrm{𝟏𝟎}}+\omega _u𝐘_{\mathrm{𝟏𝟐𝟔}}+(\stackrel{~}{\upsilon }_u+\stackrel{~}{\omega }_u)𝐘_{\mathrm{𝟏𝟐𝟎}}`$ (2.37)
$`𝐌_\nu ^{\left(\text{Dir}\right)}`$ $`=`$ $`\upsilon _u𝐘_{\mathrm{𝟏𝟎}}3\omega _u𝐘_{\mathrm{𝟏𝟐𝟔}}+(\stackrel{~}{\upsilon }_u3\stackrel{~}{\omega }_u)𝐘_{\mathrm{𝟏𝟐𝟎}}`$ (2.38)
$`𝐌_{\nu R}^{\left(\text{Maj}\right)}`$ $``$ $`M_I𝐘_{\mathrm{𝟏𝟐𝟔}}`$ (2.39)
$`𝐌_{\nu L}^{\left(\text{Maj}\right)}`$ $``$ $`{\displaystyle \frac{\omega _u^2}{M_I}}𝐘_{\mathrm{𝟏𝟐𝟔}}`$ (2.40)
Die Beziehungen (2.35-2.40) sind in $`SO(10)`$-Theorien mit intermediärer $`G_{\text{PS}}[D]`$-Symmetrie im gesamten Bereich oberhalb von $`M_I`$ gültig, während sie für die anderen $`G_I`$ in Tabelle 2.2 nur oberhalb von $`M_U`$ gelten.
In (2.35-2.40) ist zu berücksichtigen, daß $`𝐘_{\mathrm{𝟏𝟎}}`$ und $`𝐘_{\mathrm{𝟏𝟐𝟔}}`$ symmetrisch sind, während $`𝐘_{\mathrm{𝟏𝟐𝟎}}`$ antisymmetrisch ist. Um Modelle mit asymmetrischen Massenmatrizen konstruieren zu können, muß also mindestens eine 120 an der Massenerzeugung beteiligt sein. Die Faktoren $`(3)`$ sind Clebsch-Gordan-Koeffizienten, welche auf der nichttrivialen $`SU(4)_C`$-Struktur der (15,2,2) beruhen.
Wird im einfachsten Fall lediglich eine (komplexe) 10 verwendet, gilt zwischen den Massenmatrize die Beziehung
$$𝐌_e=𝐌_d𝐌_u=𝐌_\nu ^{\left(\text{Dir}\right)},$$
(2.41)
in der die $`SU(5)`$-Vorhersage $`𝐌_e=𝐌_d`$ enthalten ist. Genauso wie dort führt (2.41) aber zu falschen Werten für die Fermionmassen bei $`M_Z`$. Verwendet man Modelle mit komplizierterer Higgs-Struktur, kann man Massen und Mischungen der Fermionen korrekt beschreiben.
Die beiden Majorana-Massenmatrizen der Neutrinos sind nur bis auf konstante Vorfaktoren bekannt, die wiederum von der expliziten Form des Higgs-Potentials und den Werten der Parameter in diesem abhängen. Da es für Modelle mit nichtminimalem Higgs-Inhalt praktisch unmöglich ist, das zugehörige Potential im Detail zu analysieren, kann man lediglich plausible Annahmen über die Vorfaktoren machen. Wir werden im folgenden davon ausgehen, daß sie betragsmäßig zwischen $`1/100`$ und $`100`$ liegen. Majorana-Massenmatrizen sind aufgrund der Spinor-Struktur (1.10-1.11) der entsprechenden Massenterme stets symmetrisch. Im Rahmen von $`SO(10)`$-Modellen können sie nur durch Kopplungen an 126-Darstellungen erzeugt werden.
Die Majorana-Massen der linkshändigen Neutrinos, die durch einen VEW der $`(\overline{\mathrm{𝟏𝟎}}\mathbf{,3,1})`$ in (2.33) entstehen können, müssen sehr klein sein, da sie sonst zu beobachtbaren Effekten in durch neutrale schwache Ströme vermittelten Prozessen führen würden. Die Konstruktion eines Potentials, in dem dieser VEW verschwindet, stößt jedoch auf formale Probleme, denn in höheren Ordnungen treten Divergenzen auf, die nur dann absorbiert werden können, wenn der VEW der $`(\overline{\mathrm{𝟏𝟎}}\mathbf{,3,1})`$ ungleich Null ist . Es kann allerdings gezeigt werden, daß er von der Größenordnung $`\omega _u^2/M_I`$ sein muß und somit gegenüber den Dirac-Massen um einen Faktor $`\omega _u/M_I`$ unterdrückt ist und vernachlässigt werden kann.
Bei Skalen unterhalb von $`M_I`$ ist das Teilchenspektrum das des SM. Das Higgs-Doublett im SM ist eine bestimmte Linearkombination aus den $`SU(2)_L`$-Doubletts in den $`(\mathbf{1,2,2})`$\- und $`(\mathbf{15,2,2})`$-Darstellungen, welche an der Massenerzeugung mitwirken. Die übrigen Linearkombinationen haben Massen $`M_I`$ und treten im SM deshalb nicht in Erscheinung .
#### 2.3.4 See-Saw-Mechanismus
Die Massen der linkshändigen Neutrinos müssen, wenn sie von Null verschieden sind, sehr viel kleiner als die der geladenen Fermionen sein. Die Obergrenzen aus den Experimenten zur direkten Messung von Neutrinomassen liegen bei $`m_{\nu _e}<15`$ eV (Untersuchung des $`\beta `$-Spektrums von Tritium), $`m_{\nu _\mu }<0.17`$ MeV (aus $`\pi `$-Zerfällen) und $`m_{\nu _\tau }<18.2`$ MeV (aus $`\tau `$-Zerfällen) . Kosmologische Argumente schränken die Grenzen der Massenwerte weiter ein ; damit die kritische Dichte des Universums nicht überschritten wird, muß für die Summe der Massen leichter stabiler Neutrinos $`_im_{\nu _i}<30`$ eV gelten. Rechtshändige Neutrinos dagegen müssen sehr massiv sein, da sie experimentell nicht beobachtet werden.
Der See-Saw-Mechanismus liefert in Modellen mit rechtshändigen Neutrinos und einer Skala $`M_IM_Z`$, bei der die Leptonzahlerhaltung verletzt wird, eine natürliche Erklärung für die oben angegebenen Eigenschaften des Neutrinomassenspektrums. Seine Wirkungsweise soll zunächst am einfachen Fall einer Fermiongeneration verdeutlicht werden.
Der allgemeinste Massenterm für das Neutrino $`\nu `$ lautet:
$`_M`$ $`=`$ $`{\displaystyle \frac{1}{2}}m_R^{\left(\text{Maj}\right)}\overline{\nu }_L^C\nu _R{\displaystyle \frac{1}{2}}m_L^{\left(\text{Maj}\right)}\overline{\nu }_L\nu _R^Cm^{\left(\text{Dir}\right)}\overline{\nu }_L\nu _R+\text{h.c.}`$ (2.42)
$`=`$ $`{\displaystyle \frac{1}{2}}\left(\begin{array}{cc}\overline{\nu }_L& \overline{\nu }_L^C\end{array}\right)𝐌\left(\begin{array}{c}\nu _R^C\\ \nu _R\end{array}\right)+\text{h.c.}`$
mit
$$𝐌=\left(\begin{array}{cc}m_L^{\left(\text{Maj}\right)}& m^{\left(\text{Dir}\right)}\\ m^{\left(\text{Dir}\right)}& m_R^{\left(\text{Maj}\right)}\end{array}\right)$$
(2.43)
In $`SO(10)`$-GUTs mit intermediärer Symmetrie gelten die Relationen
$$\frac{m^{\left(\text{Dir}\right)}}{m_R^{\left(\text{Maj}\right)}}\mathrm{\hspace{0.17em}1}\text{und}\frac{m_L^{\left(\text{Maj}\right)}}{m^{\left(\text{Dir}\right)}}\mathrm{\hspace{0.17em}1},$$
(2.44)
da, wie im vorigen Abschnitt erläutert wurde, $`m^{\left(\text{Dir}\right)}`$ dieselbe Größenordnung wie die Masse des positiv geladenen Quarks hat und $`m_R^{\left(\text{Maj}\right)}`$ von der Ordnung der $`BL`$-brechenden Skala $`M_I`$ ist; $`m_L^{\left(\text{Maj}\right)}`$ wiederum ist gegenüber $`m^{\left(\text{Dir}\right)}`$ stark unterdrückt.
Die Diagonalisierung von M führt (abgesehen von Korrekturen $`m^{\left(\text{Dir}\right)}/m_R^{\left(\text{Maj}\right)}`$) auf die Masseneigenzustände
$$\nu _1=\frac{1}{\sqrt{2}}(\nu _L+\nu _R^C)\text{und}\nu _2=\frac{1}{\sqrt{2}}(\nu _R+\nu _L^C)$$
(2.45)
und ihre Massen
$`m_{\nu _1}`$ $``$ $`m_L^{\left(\text{Maj}\right)}\left({\displaystyle \frac{m^{\left(\text{Dir}\right)}}{m_R^{\left(\text{Maj}\right)}}}\right)m^{\left(\text{Dir}\right)}\left({\displaystyle \frac{m^{\left(\text{Dir}\right)}}{m_R^{\left(\text{Maj}\right)}}}\right)m^{\left(\text{Dir}\right)}`$ (2.46)
$`m_{\nu _2}`$ $``$ $`m_R^{\left(\text{Maj}\right)}+\left({\displaystyle \frac{m^{\left(\text{Dir}\right)}}{m_R^{\left(\text{Maj}\right)}}}\right)m^{\left(\text{Dir}\right)}m_R^{\left(\text{Maj}\right)}`$ (2.47)
Anstelle eines Dirac-Neutrinos $`\nu =\nu _L+\nu _R`$ erhält man auf diese Weise zwei Majorana-Neutrinos $`\nu _1`$ und $`\nu _2`$ (es gilt $`\nu _{\mathrm{1,2}}^C=\nu _{\mathrm{1,2}}`$) mit Massen $`m_{\nu _{\mathrm{1,2}}}`$, wobei $`m_{\nu _1}M_Zm_{\nu _2}`$ ist. Das leichte Majorana-Neutrino $`\nu _1`$ entspricht dem Neutrino des SM; das rechtshändige Antineutrino des SM ist identisch mit $`\nu _{1R}\nu _{1R}^C`$. Das schwere Neutrino $`\nu _2`$ ist aufgrund seiner Masse $`m_{\nu _2}M_I`$ nicht beobachtbar.
Aus (2.46) und (2.47) folgt die „See-Saw“-Relation
$$m_{\nu _1}m_{\nu _2}m^{\left(\text{Dir}\right)}m^{\left(\text{Dir}\right)}M_Z^2;$$
(2.48)
je schwerer $`\nu _2`$ ist, desto leichter wird $`\nu _1`$.
Die Verallgemeinerung von (2.42) auf den Fall von $`n`$ Fermionfamilien liefert
$$_M=\frac{1}{2}\left(\begin{array}{cc}\overline{\nu }_L& \overline{\nu }_L^C\end{array}\right)𝐌\left(\begin{array}{c}\nu _R^C\\ \nu _R\end{array}\right)+\text{h.c.}$$
(2.49)
mit $`n`$-dimensionalen Vektoren $`\nu _{L,R}`$ und der symmetrischen ($`2n\times 2n`$)-Matrix
$$𝐌=\left(\begin{array}{cc}𝐌_{\nu L}^{\left(\text{Maj}\right)}& 𝐌_\nu ^{\left(\text{Dir}\right)}\\ \left(𝐌_\nu ^{\left(\text{Dir}\right)}\right)^T& 𝐌_{\nu R}^{\left(\text{Maj}\right)}\end{array}\right)$$
(2.50)
Hier sind die nichtverschwindenden Einträge von $`𝐌_{\nu L}^{\left(\text{Maj}\right)}`$ betragsmäßig sehr viel kleiner als die von $`𝐌_\nu ^{\left(\text{Dir}\right)}`$, und diese sind wiederum sehr viel kleiner als die von $`𝐌_{\nu R}^{\left(\text{Maj}\right)}`$ (siehe auch (2.38-2.40)). Indem man M auf blockdiagonale Form bringt, erhält man $`n`$ leichte Majorana-Neutrinos $`\nu _i`$ mit der symmetrischen Massenmatrix
$$𝐌_\nu 𝐌_{\nu L}^{\left(\text{Maj}\right)}𝐌_\nu ^{\left(\text{Dir}\right)}\left(𝐌_{\nu R}^{\left(\text{Maj}\right)}\right)^1\left(𝐌_\nu ^{\left(\text{Dir}\right)}\right)^T𝐌_\nu ^{\left(\text{Dir}\right)}\left(𝐌_{\nu R}^{\left(\text{Maj}\right)}\right)^1\left(𝐌_\nu ^{\left(\text{Dir}\right)}\right)^T,$$
(2.51)
der sogenannten See-Saw-Matrix, und $`n`$ schwere Majorana-Neutrinos $`N_i`$ mit der Massenmatrix
$$𝐌_N𝐌_{\nu R}^{\left(\text{Maj}\right)}+𝐌_\nu ^{\left(\text{Dir}\right)}\left(𝐌_{\nu R}^{\left(\text{Maj}\right)}\right)^1\left(𝐌_\nu ^{\left(\text{Dir}\right)}\right)^T𝐌_{\nu R}^{\left(\text{Maj}\right)}$$
(2.52)
Die Diagonalisierung von $`𝐌_\nu `$ und $`𝐌_N`$
$$𝐍_L^{}𝐌_\nu 𝐍_L=𝐌_\nu ^{\left(D\right)},𝐍_R^{}𝐌_N𝐍_R=𝐌_N^{\left(D\right)}$$
(2.53)
liefert für $`n=3`$ die physikalischen Neutrinos $`\nu _e`$, $`\nu _\mu `$, $`\nu _\tau `$ beziehungsweise $`N_e`$, $`N_\mu `$, $`N_\tau `$ sowie deren Massen und Mischungen. Letztere haben nun Auswirkungen auf die leptonischen Anteile der geladenen schwachen Ströme im SM, welche sich beim Übergang von Wechselwirkungs- zu Masseneigenzuständen analog zu (1.31) gemäß
$`\overline{\nu }_L^{a\left(0\right)}\gamma ^\mu e_L^{a\left(0\right)}W_\mu ^+`$ $``$ $`\left(𝐍_L^{}𝐄_L\right)_{ab}\overline{\nu }_L^b\gamma ^\mu e_L^aW_\mu ^+`$ (2.54)
$`\overline{e}_L^{a\left(0\right)}\gamma ^\mu \nu _L^{a\left(0\right)}W_\mu ^{}`$ $``$ $`\overline{e}_L^a\gamma ^\mu \left(𝐄_L^{}𝐍_L\right)_{ab}\nu _L^bW_\mu ^{}`$ (2.55)
ändern ($`a`$ und $`b`$ sind Familienindizes). Die unitäre Matrix $`𝐔𝐄_L^{}𝐍_L`$ entspricht der CKM-Matrix im Quarksektor und kann in Neutrino-Oszillationsexperimenten gemessen werden. Die beobachtbaren Effekte der Neutrinomassen und -mischungen werden im nächsten Kapitel erläutert.
#### 2.3.5 Zusammenfassung
Im Vergleich zur $`SU(5)`$-GUT sind $`SO(10)`$-Modelle mit intermediärer Symmetrie phänomenologisch sehr erfolgreich, und es sind bis heute keine experimentellen Resultate bekannt, die sie ausschließen. Zusätzlich zu den positiven Eigenschaften, welche schon die $`SU(5)`$ besaß, kommen hier folgende Vorzüge:
* Alle Fermionen einer Generation, einschließlich des rechtshändigen Neutrinos, liegen in einer irreduziblen Darstellung.
* Die Eigenschaften der Neutrinos, insbesondere deren Massen und Mischungen, können durch den See-Saw-Mechanismus befriedigend erklärt werden.
* Die $`SO(10)`$ ist automatisch anomaliefrei, enthält $`BL`$ als lokale Symmetrie und besitzt rechts-links-symmetrische Untergruppen.
Die Probleme, welche weder in $`SO(10)`$-Modellen noch in anderen GUTs gelöst werden können, sind im letzten Abschnitt zusammengefaßt.
### 2.4 $`E_6`$
Die $`E_6`$ ist die einzige exzeptionelle Lie-Gruppe mit komplexen Darstellungen. Zu den maximalen Untergruppen der $`E_6`$ gehören die für die Symmetriebrechung nach $`G_{\text{SM}}`$ in Frage kommenden $`SO(10)U(1)`$ und $`[SU(3)]^3`$. Letztere kann man als $`SU(3)_LSU(3)_RSU(3)_C`$-Symmetrie interpretieren.
Die Fermionen einer Generation liegen in der komplexen Fundamentaldarstellung 27, welche sich bezüglich der $`SO(10)`$ und $`[SU(3)]^3`$ gemäß
$`\mathrm{𝟐𝟕}`$ $``$ $`\mathrm{𝟏𝟔}\mathrm{𝟏𝟎}\mathrm{𝟏}`$ (2.56)
$`\mathrm{𝟐𝟕}`$ $``$ $`(\overline{\mathrm{𝟑}}\mathbf{,3,1})(\mathbf{3,1,3})(\mathrm{𝟏},\overline{\mathrm{𝟑}},\overline{\mathrm{𝟑}})`$ (2.57)
verzweigt. Im Gegensatz zu den $`SU(5)`$\- und $`SO(10)`$-Modellen enthält die $`E_6`$-GUT zusätzliche Fermionen, die im SM nicht beobachtet werden und somit superschwer sein müssen. In der 10 liegen ein Quark $`D`$ mit Ladung $`1/3`$, ein leptonisches $`SU(2)_L`$-Doublett $`(N,E)`$ sowie die zugehörigen Antiteilchen; das $`SO(10)`$-Singulett wird mit $`L`$ bezeichnet.
Die adjungierte Darstellung 78 der $`E_6`$ enthält die Eichbosonen der Theorie. Die im Vergleich zur $`SO(10)`$ neu hinzukommenden Eichbosonen koppeln sowohl an die superschweren als auch an die SM-Fermionen und liefern zu Nukleonzerfällen keine Beiträge.
Die Fermionmassen besitzen das $`E_6`$-Transformationsverhalten
$$\mathrm{𝟐𝟕}\mathrm{𝟐𝟕}=\mathbf{\hspace{0.33em}27}\mathrm{𝟑𝟓𝟏}\mathrm{𝟑𝟓𝟏}^{},$$
(2.58)
wobei die 351 und die $`\mathrm{𝟑𝟓𝟏}^{}`$ inäquivalent sind. Unter $`SO(10)`$ verzweigen sich diese beiden Darstellungen wie
$`\mathrm{𝟑𝟓𝟏}`$ $``$ $`\mathrm{𝟏}\mathrm{𝟏𝟎}\mathrm{𝟏𝟐𝟔}\mathrm{}\mathrm{}`$ (2.59)
$`\mathrm{𝟑𝟓𝟏}^{}`$ $``$ $`\mathrm{𝟏𝟎}\mathrm{𝟏𝟐𝟎}\mathrm{}\mathrm{}`$ (2.60)
Die 351 kann $`E_6`$ in die $`SO(10)`$ brechen, ohne daß die SM-Fermionen Massen der Größenordnung $`M_U`$ bekommen; in diesem Brechungsschritt werden nur Massen für die neuen Fermionen erzeugt. Die Brechung $`G_{\text{SM}}SU(3)_CU(1)_{\text{em}}`$ und die Erzeugung der SM-Massen werden durch die Higgs-Darstellung $`\overline{\mathrm{𝟐𝟕}}`$ realisiert. Um asymmetrische Massenmatrizen erhalten zu können, muß eine $`\mathrm{𝟑𝟓𝟏}^{}`$ vorhanden sein.
Die kleinste Darstellung, welche die Symmetriebrechung $`E_6[SU(3)]^3`$ realisieren kann, ist die 650. In beiden Brechungsschemata ergeben sich Lebensdauern für das Proton, die weit oberhalb des experimentell zugänglichen Bereichs liegen.
Offensichtlich ist die allgemeine Struktur von $`E_6`$-GUTs wesentlich komplizierter als die von $`SO(10)`$-Modellen; insbesondere die Existenz exotischer superschwerer Fermionen kommt neu hinzu. Dem steht als Vorteil die Möglichkeit gegenüber, Fermionmassen und -mischungen durch Strahlungskorrekturen zu erzeugen .
### 2.5 Probleme und Grenzen von GUTs
Wie in diesem Kapitel dargestellt wurde, können GUTs viele Schwächen des SM erfolgreich beheben. Ferner liefern sie eine Reihe von zumindest grundsätzlich überprüfbaren Vorhersagen wie zum Beispiel die Instabilität der Nukleonen oder die Neutrinomassen und -mischungen. Es gibt jedoch auch Probleme, die im Rahmen von nichtsupersymmetrischen GUT-Modellen nicht gelöst werden können. Man kann sie in rein technische und fundamentale Probleme unterteilen.
Zu den ersteren gehören das schon erwähnte Hierarchieproblem und die daraus resultierenden Feineinstellungen verschiedener Parameter. Diese Feineinstellungen sind bezüglich Strahlungskorrekturen instabil, so daß sie für jede Ordnung der Störungstheorie durchgeführt werden müssen. Das $`\theta _{\text{QCD}}`$-Problem ($`\theta _{\text{QCD}}10^9`$) gehört ebenfalls in diese Kategorie.
Ungelöste fundamentale Fragestellungen sind die Abwesenheit der Gravitationswechselwirkung in GUTs, eine fehlende Erklärung für die Existenz dreier Fermionfamilien und ein hohes Maß an Unbestimmtheit im Higgs-Sektor der Modelle. Auch der Ursprung der Massenhierarchie der Fermionen kann nicht geklärt werden, obwohl es in GUTs im Gegensatz zum SM Beziehungen zwischen den Massenmatrizen der verschiedenen Fermionarten gibt.
Eine wirklich befriedigende Lösung dieser Probleme im Rahmen einer fundamentaleren Theorie als den hier geschilderten GUT-Modellen steht jedoch bis heute aus.
## Kapitel 3 Neutrino-Oszillationen
### 3.1 Theoretische Grundlagen
Neutrino-Oszillationen können auftreten, wenn die Masseneigenzustände der Neutrinos nicht mit den Wechselwirkungseigenzuständen übereinstimmen. Dies ist in $`SO(10)`$-Modellen mit intermediärer Symmetrie im allgemeinen der Fall, da der See-Saw-Mechanismus aus Abschnitt 2.3.4 für die leichten Majorana-Neutrinos $`\nu _\alpha `$ ($`\alpha =e,\mu ,\tau `$) die üblicherweise nichtdiagonale Massenmatrix (2.51) liefert. Wenn die Mischungsmatrix $`𝐍_L\mathrm{𝟏}`$ ist, sind die $`\nu _\alpha `$ Linearkombinationen der Neutrinos $`\nu _k`$ ($`k`$=1,2,3) mit den Massen $`m_k`$, wobei $`m_1m_2m_3`$ gelten soll:
$$\nu _{\alpha L,R}=\underset{k=1}{\overset{3}{}}𝐔_{\alpha k}\nu _{kL,R},𝐔𝐄_L^{}𝐍_L$$
(3.1)
Die schweren Neutrinos $`N_\alpha `$ entkoppeln bei Skalen $`\mu M_I`$ und spielen für die in diesem Kapitel geschilderten Phänomene keine Rolle.
#### 3.1.1 Vakuumoszillationen
Zunächst sollen Oszillationen im Vakuum betrachtet werden. Der quantenmechanische Zustand eines zur Zeit $`t=0`$ in einem elektroschwachen Prozeß erzeugten relativistischen Neutrinos $`\nu _\alpha `$ mit Impuls $`pm_k`$ ist durch
$$|\nu _\alpha (x,t=0)=\underset{k}{}𝐔_{\alpha k}e^{ipx}|\nu _k(x,t=0)$$
(3.2)
gegeben. Die Zeitentwicklung der Zustände $`|\nu _k(x,t=0)`$ gemäß der Schrödinger-Gleichung liefert unter Verwendung von
$$E=E(\nu _k)=\sqrt{p^2+m_k^2}p+\frac{m_k^2}{2p}\text{und}xct(c1)$$
(3.3)
für $`|\nu _\alpha (x)`$ die Beziehung
$$|\nu _\alpha (x)=\underset{k}{}𝐔_{\alpha k}e^{i(m_k^2/2E)x}|\nu _k(0)$$
(3.4)
Drückt man $`|\nu _k(x=0)`$ als Linearkombination der $`|\nu _\alpha (x=0)`$ aus, so ergibt sich
$$|\nu _\alpha (x)=\underset{k,\beta }{}𝐔_{\alpha k}𝐔_{\beta k}^{}e^{i(m_k^2/2E)x}|\nu _\beta (0)$$
(3.5)
Mit $`P(\nu _\alpha \nu _\beta ,L)`$ sei die Wahrscheinlichkeit bezeichnet, daß ein Neutrino der Energie $`E`$, welches bei $`x=0`$ als $`\nu _\alpha `$ erzeugt wurde, in einer Entfernung $`L`$ von der Quelle als $`\nu _\beta `$ detektiert wird. Dann gilt unter Berücksichtigung der Unitarität von U (die in Modellen mit See-Saw-Mechanismus bis auf Korrekturen der Ordnung $`M_Z/M_I`$ gegeben ist):
$$P(\nu _\alpha \nu _\beta ,L)=|\nu _\beta (L)|\nu _\alpha (0)|^2=|\delta _{\alpha \beta }+\underset{k=2}{\overset{3}{}}𝐔_{\alpha k}𝐔_{\beta k}^{}[e^{i(\mathrm{\Delta }m_{k1}^2/2E)L}1]|^2$$
(3.6)
Hierbei ist $`\mathrm{\Delta }m_{k1}^2m_k^2m_1^2`$. Aus der $`CPT`$-Invarianz folgen ferner die Beziehungen
$$P(\nu _\alpha \nu _\beta ,L)=P(\overline{\nu }_\beta \overline{\nu }_\alpha ,L)\text{und}P(\nu _\alpha \nu _\alpha ,L)=P(\overline{\nu }_\alpha \overline{\nu }_\alpha ,L)$$
(3.7)
Die Übergangswahrscheinlichkeit für $`\nu _\alpha \nu _\beta `$ hängt demnach von den Elementen von U, von zwei unabhängigen Differenzen der Massenquadrate und vom Parameter $`L/E`$ ab. Sie kann dann von Null verschieden sein, wenn $`𝐔\mathrm{𝟏}`$ ist und für mindestens ein $`k`$ $`\mathrm{\Delta }m_{k1}^2E/L`$ gilt. Ist dies der Fall, so wird ein Neutrinostrahl, welcher bei $`x=0`$ nur aus $`\nu _\alpha `$ besteht, an einem Ort $`x=L`$ auch Neutrinos der Art $`\nu _{\beta \alpha }`$ beinhalten. Ein nur für den Nachweis von $`\nu _\alpha `$ geeignetes Experiment wird dann ein Neutrinodefizit feststellen.
Die Bedeutung von (3.6) wird besonders klar, wenn man den Fall zweier Neutrino-Arten betrachtet. Dann besitzt U die einfache Parametrisierung
$$𝐔=\left(\begin{array}{cc}\mathrm{cos}\theta & \mathrm{sin}\theta \\ \mathrm{sin}\theta & \mathrm{cos}\theta \end{array}\right)$$
(3.8)
(U kann zwar $`CP`$-verletzende Phasen enthalten, die aber auf Neutrino-Oszillationen keinen Einfluß haben) und es gilt:
$`P(\nu _\alpha \nu _\beta ,L)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{sin}^22\theta \left(\mathrm{\hspace{0.17em}1}\mathrm{cos}(\mathrm{\Delta }m^2L/2E)\right)(\alpha \beta )`$ (3.9)
$`P(\nu _\alpha \nu _\alpha ,L)`$ $`=`$ $`P(\nu _\beta \nu _\beta ,L)=\mathrm{\hspace{0.33em}1}P(\nu _\alpha \nu _\beta ,L)`$ (3.10)
Die Übergangswahrscheinlichkeit (3.9) ist eine periodische Funktion von $`L/E`$; dieses Phänomen wird als Neutrino-Oszillation bezeichnet. Die Amplitude der Oszillation ist gleich $`\mathrm{sin}^22\theta `$, hängt also allgemein von den Einträgen von U ab, und die Oszillationslänge ist durch $`L^{\text{osc}}=4\pi E/\mathrm{\Delta }m^2`$ gegeben. Oszillationseffekte können demnach beobachtet werden, sofern $`LL^{\text{osc}}`$ ist.
Bei der Datenanalyse von Oszillationsexperimenten wird im allgemeinen von nur zwei beteiligten Neutrino-Arten ausgegangen und (3.9-3.10) verwendet. Die graphische Auswertung erfolgt dann in $`\mathrm{sin}^22\theta \mathrm{\Delta }m^2`$-Diagrammen, wobei die Nichtbeobachtung von Oszillationen abhängig von der Nachweisempfindlichkeit des Experiments bestimmte Bereiche des Parameterraums ausschließt, während positive Resultate zu mehr oder weniger eng begrenzten erlaubten Parameterbereichen führen.
Oszillationsexperimente können lediglich Aussagen über die Größen $`|𝐔|`$ und $`\mathrm{\Delta }m_{jk}^2`$ machen. Auf diese Weise lassen sich die Verhältnisse je zweier Neutrinomassen bestimmen, nicht aber die Massenwerte selbst.
#### 3.1.2 Oszillationen in Materie
Wenn Neutrinos Materie durchqueren, kann der sogenannte Mikheyev-Smirnov-Wolfenstein-Effekt (MSW-Effekt) auftreten, welcher das Oszillationsverhalten im Vergleich zum Vakuumfall modifiziert. Ursache des MSW-Effekts ist die Tatsache, daß in Materie für die Teilchendichten $`N_i`$ der geladenen Leptonen $`N_eN_{\mu ,\tau }0`$ gilt. Während alle Neutrinos die gleiche Wechselwirkung über neutrale schwache Ströme erfahren, wechselwirken nur die Elektron-Neutrinos über geladene schwache Ströme mit der Materie. Das führt zu einer Flavour-Asymmetrie in der Vorwärts-Streuung von Neutrinos, die eine Phasenverschiebung zur Folge hat und sich somit auf die Zeitentwicklung des Gesamtsystems auswirkt .
Eine formale Analyse dieses Sachverhalts liefert für zwei Neutrino-Arten wieder die Beziehungen (3.9-3.10), wobei allerdings die Vakuumgrößen $`\mathrm{sin}^22\theta `$ und $`L^{\text{osc}}`$ durch die neuen Größen
$`\mathrm{sin}^22\theta _M`$ $`=`$ $`{\displaystyle \frac{\mathrm{sin}^22\theta }{(\mathrm{cos}2\theta L^{\text{osc}}/L_0)^2+\mathrm{sin}^22\theta }}\text{und}`$ (3.11)
$`L_M^{\text{osc}}`$ $`=`$ $`{\displaystyle \frac{L^{\text{osc}}}{\sqrt{(\mathrm{cos}2\theta L^{\text{osc}}/L_0)^2+\mathrm{sin}^22\theta }}}`$ (3.12)
mit $`L_0=2\pi /(\sqrt{2}G_FN_e)`$ ersetzt werden müssen. Daraus ergeben sich folgende Spezialfälle:
* $`L^{\text{osc}}L_0`$ (geringe Elektronendichte):
$$\mathrm{sin}^22\theta _M\mathrm{sin}^22\theta ,L_M^{\text{osc}}L^{\text{osc}}$$
(3.13)
Die Materie-Effekte sind vernachlässigbar klein.
* $`L^{\text{osc}}L_0`$ (hohe Elektronendichte):
$$\mathrm{sin}^22\theta _M(L_0/L^{\text{osc}})^2\mathrm{sin}^22\theta \mathrm{\hspace{0.33em}0},L_M^{\text{osc}}L_0$$
(3.14)
Die Vakuumparameter sind stark unterdrückt; die Amplitude ist verschwindend klein und die Oszillationslänge ergibt sich unabhängig von $`L^{\text{osc}}`$.
* $`L^{\text{osc}}/L_0=\mathrm{cos}2\theta `$:
$$\mathrm{sin}^22\theta _M=\mathrm{\hspace{0.33em}1},L_M^{\text{osc}}=L^{\text{osc}}/\mathrm{sin}2\theta $$
(3.15)
Die Oszillationsamplitude ist maximal, auch wenn der entsprechende Vakuumwert $`\mathrm{sin}^22\theta `$ sehr klein ist. Besitzt die Materieverteilung eine räumliche Ausdehnung der Größenordnung $`L_M^{\text{osc}}`$, so kann die Übergangswahrscheinlichkeit durchaus $`1`$ sein. Dieses Resonanzverhalten von (3.11) wird als MSW-Effekt bezeichnet .
Die Materie-Effekte auf das Oszillationsverhalten sind besonders für den Fall der im Inneren der Sonne erzeugten Elektron-Neutrinos von Bedeutung.
### 3.2 Experimentelle Situation
Die Neutrino-Oszillationsexperimente lassen sich grob in zwei Kategorien unterteilen. Die „Appearence“-Experimente untersuchen das Erscheinen von $`\nu _{\beta \alpha }`$ in einem Neutrinostrahl, welcher bei seiner Entstehung ausschließlich aus $`\nu _\alpha `$-Neutrinos bestand, während „Disappearence“-Experimente nach einer Verringerung des $`\nu _\alpha `$-Flusses in einem solchen Strahl suchen. Weiterhin kann man die Experimente nach den Neutrinoquellen unterscheiden; es finden Neutrinos aus Beschleunigern, Kernreaktoren und natürlichen Quellen Verwendung. Ausführliche Referenzen zu den hier erwähnten Oszillationsexperimenten finden sich in . Bis heute sind drei Indizien bekannt, die auf Neutrino-Oszillationen hindeuten.
#### 3.2.1 Sonnen-Neutrinos
Struktur und Dynamik der Sonne werden durch das Standard-Sonnenmodell (SSM) beschrieben und gelten als gut verstanden. Die Energie wird im Kern der Sonne über die thermonuklearen $`pp`$\- und CNO-Reaktionsketten erzeugt; der zugrundeliegende Fusionsprozeß ist in beiden Fällen
$$4p\alpha +2e^++2\nu _e+26.7\text{MeV}$$
(3.16)
Das Energiespektrum der Elektron-Neutrinos ist inhomogen und reicht je nach erzeugender Reaktion von 100 keV bis zu 15 MeV. Da die Unsicherheiten in den SSM-Vorhersagen über den Neutrinofluß relativ klein sind, eignet sich dieser gut für die Suche nach Oszillationen $`\nu _e\nu _{\mu ,\tau }`$. Alle bisher durchgeführten Messungen haben beträchtliche Differenzen zwischen tatsächlich gemessenen Ereignissen $`N_{\text{Exp}}`$ und den aufgrund von Monte-Carlo-Simulationen erwarteten Ereignissen $`N_{\text{SSM}}`$ festgestellt; die Resultate sind in Tabelle 3.1 dargestellt .
Homestake verwendet zum Nachweis 600 t $`\text{C}_2\text{Cl}_4`$, GALLEX und SAGE bestehen aus 30 t Galliumchlorid beziehungsweise 60 t metallischem Gallium, Kamiokande und das Nachfolge-Experiment Super-Kamiokande sind Wasser-Čerenkovdetektoren.
Die Analyse der Daten führt zu dem Ergebnis, daß sowohl Vakuumoszillationen auf dem Weg zwischen der Sonne und der Erde als auch der MSW-Effekt im Sonneninneren das Neutrinodefizit erklären können;
es kommen $`\nu _e\nu _\mu `$\- und $`\nu _e\nu _\tau `$-Übergänge in Frage . Die Abbildungen 3.1(a) und 3.1(b) geben die erlaubten Bereiche (mit einem Confidence Level (CL) von 99%) im $`\mathrm{sin}^22\theta \mathrm{\Delta }m^2`$-Parameterraum an, wobei die Punkte den statistisch wahrscheinlichsten Lösungen entsprechen:
* eine MSW-Lösung mit großem Mischungswinkel bei $`\mathrm{sin}^22\theta 0.76`$ und $`\mathrm{\Delta }m^2210^5\text{eV}^2`$.
* eine MSW-Lösung mit kleinem Mischungswinkel bei $`\mathrm{sin}^22\theta 610^3`$ und $`\mathrm{\Delta }m^2510^6\text{eV}^2`$.
* eine MSW-Lösung mit kleinem $`\mathrm{\Delta }m^2`$ bei $`\mathrm{sin}^22\theta 0.96`$ und $`\mathrm{\Delta }m^2810^8\text{eV}^2`$.
* eine Vakuumoszillations-Lösung bei $`\mathrm{sin}^22\theta 0.75`$ und $`\mathrm{\Delta }m^2810^{11}\text{eV}^2`$.
Die MSW-Lösung mit kleinem $`\mathrm{\Delta }m^2`$ ist in ihrer Statistik wesentlich unwahrscheinlicher als die anderen, da sie für 95 % CL nicht mehr akzeptabel ist. Damit die Vakuumoszillations-Lösung das Neutrinodefizit erklären kann, ist eine nicht sehr natürlich wirkende Feineinstellung zwischen der Oszillationslänge und dem Sonne-Erde-Abstand erforderlich. Deswegen werden in dieser Arbeit nur die ersten beiden Lösungen in Betracht gezogen.
Weitere Experimente zur Untersuchung des Sonnen-Neutrinodefizits sind in Vorbereitung. Der Nachfolger von GALLEX ist das unterirdische Gallium Neutrino Observatory (GNO) im Gran Sasso-Labor (Italien) mit angestrebten 100 t Ga im Jahre 2002. Borexino ist ein organischer Flüssig-Szintillator mit einer Empfindlichkeit $`E_\nu >0.25`$ MeV und wird ebenfalls in Gran Sasso untergebracht. Das Sudbury Neutrino Observatory (SNO) in Kanada wird ein Schwerwasser-Čerenkovdetektor aus 1000 t D<sub>2</sub>O sein, welcher für $`E_\nu >5`$ MeV sensibel ist.
#### 3.2.2 Atmosphärische Neutrinos
Treffen hochenergetische Protonen der kosmischen Strahlung auf die obere Erdatmosphäre, so kollidieren sie mit den Kernen der Luftmoleküle. In diesen hadronischen Stoßprozessen entstehen zahlreiche Pionen und Kaonen, die gemäß
$`\pi ^+,K^+\mu ^++\nu _\mu `$ $`\pi ^{},K^{}\mu ^{}+\overline{\nu }_\mu `$
$`\text{ }e^++\nu _e+\overline{\nu }_\mu `$ $`\text{ }e^{}+\overline{\nu }_e+\nu _\mu `$ (3.17)
zerfallen; die dabei erzeugten Neutrinos bezeichnet man als atmosphärische Neutrinos. Wegen (3.17) erwartet man für das zahlenmäßige Verhältnis von Myon- zu Elektron-Neutrino-Ereignissen in einem Detektor in guter Näherung
$$(\mu /e)\frac{N(\nu _\mu +\overline{\nu }_\mu )}{N(\nu _e+\overline{\nu }_e)}=\mathrm{\hspace{0.33em}2}$$
(3.18)
Obwohl der exakte Wert von $`(\mu /e)`$ von der Neutrino-Energie abhängt und von 2 abweichen kann, erlauben detaillierte Analysen mit Hilfe von Monte-Carlo-Simulationen relativ präzise Vorhersagen. Diese bieten einen weiteren Ansatzpunkt für die Suche nach Neutrino-Oszillationen und wurden vom Kamiokande-Experiment und dem größeren Nachfolger
Super-Kamiokande in Japan überprüft. Beides sind unterirdische Wasser-Čerenkovdetektoren; Kamiokande bestand aus 3 kt und Super-Kamiokande besteht aus 50 kt reinen Wassers in 1000 m Tiefe. Als Nachweisreaktionen dienen
$$\nu _e+ne^{}+p,\overline{\nu }_e+pe^++n,\nu _\mu +n\mu ^{}+p,\overline{\nu }_\mu +p\mu ^++n$$
(3.19)
Beide Experimente haben für das Verhältnis $`R`$ von gemessenem $`(\mu /e)_{\text{data}}`$ zu theoretisch erwartetem $`(\mu /e)_{\text{MC}}`$ erhebliche Abweichungen vom Wert 1 festgestellt, die auf einem Defizit von $`(\nu _\mu ,\overline{\nu }_\mu )`$-Ereignissen beruhen. Tabelle 3.2 faßt die Resultate zusammen, wobei der erste Fehler von $`R`$ statistischen und der zweite systematischen Ursprungs ist.
Eine mögliche Erklärung für diese Anomalie können $`\nu _\mu \nu _\tau `$-Oszillationen liefern. Die Oszillation in Elektron-Neutrinos ist aufgrund der Resulate des Reaktor-Experiments CHOOZ ausgeschlossen, welches nach $`\overline{\nu }_e\overline{\nu }_{\mu ,\tau }`$-Übergängen sucht, im Rahmen seines Meßbereichs ($`E_\nu 3`$ MeV, $`L=1`$ km) aber keine solchen gefunden hat.
Die Oszillationslösung wird durch die beobachtete Abhängigkeit des $`(\nu _\mu ,\overline{\nu }_\mu )`$-Defizits vom Zenit-Winkel $`\theta `$ (nicht zu verwechseln mit dem gleichnamigen Mischungswinkel) unterstützt (Abbildung 3.2(a)). Während die Neutrinos, welche den Detektor senkrecht von oben erreichen ($`\theta =0`$), eine Strecke von $`L20`$ km zurückgelegt haben, müssen die den Detektor von unten erreichenden die gesamte Erde durchqueren, das heißt $`L13000`$ km für $`\theta =\pi `$. Je größer $`\theta `$ und somit auch $`L`$ sind, desto mehr Myon-(Anti-)Neutrinos wandeln sich in Tau-(Anti-)Neutrinos um.
Abbildung 3.2(b) zeigt den erlaubten Parameterbereich für die Oszillationslösung; die wahrscheinlichste Lösung zur Erklärung der Super-Kamiokande-Resultate liegt bei $`\mathrm{sin}^22\theta =1.0`$ und $`\mathrm{\Delta }m^2=2.210^3`$ eV. Die Annahme eines hierarchischen Massenschemas $`m_1m_2m_3`$ für die Neutrinos führt zu $`m_3\sqrt{\mathrm{\Delta }m_{\text{atm}}^2}0.05`$ eV, was bedeuten würde, daß Neutrinos zur dunklen Materie keinen nennenswerten Beitrag liefern.
Unter Berücksichtigung der $`\nu _\mu \nu _\tau `$-Oszillationslösung für die Anomalie der atmosphärischen Neutrinos ist das Sonnen-Neutrinoproblem demnach auf $`\nu _e\nu _\mu `$-Oszillationen zurückzuführen, da die erlaubten Parameterbereiche in den Abbildungen 3.1 und 3.2(b) sich in keinem Fall überlappen.
Zur Ergänzung von Super-Kamiokande, das weiterhin Daten aufnimmt, sind mehrere sogenannte „Long Baseline“-Experimente geplant. Hierbei sollen Neutrinostrahlen mit wohldefinierten Eigenschaften, die mittels Stoßprozessen an Teilchenbeschleunigern erzeugt werden, in einigen hundert km Entfernung detektiert werden. Erste Resulte sind jedoch nicht vor 2001 zu erwarten.
#### 3.2.3 LSND und KARMEN
Das LSND-Beschleunigerexperiment in Los Alamos verwendet Neutrinos, die durch das Auftreffen eines relativistischen Protonenstrahls auf ein Wasser-Target entstehen. In den Stoßprozessen werden $`\pi ^+`$-Mesonen erzeugt, welche anschließend zerfallen; gemäß (3.17) enthält der resultierende Neutrinostrahl keine Elektron-Antineutrinos. Im 30 m entfernten Flüssig-Szintillationsdetektor werden allerdings Photonen nachgewiesen, welche aus der Reaktion $`\overline{\nu }_e+pe^++n`$ und nachfolgendem $`e^++e^{}\gamma `$ beziehungsweise $`n+pd+\gamma `$ stammen. Eine mögliche Erklärung hierfür können $`\overline{\nu }_\mu \overline{\nu }_e`$-Oszillationen liefern. Der erlaubte Parameterbereich ist in Abbildung 3.3 dargestellt, wobei der dunkelgraue Bereich 90 % CL und der hellgraue 99 % CL kennzeichnet.
Eine Reihe von weiteren Reaktor- und Beschleunigerexperimenten haben $`\overline{\nu }_\mu \overline{\nu }_e`$-Oszillationen untersucht, aber im Rahmen ihrer Meßbereiche keine entsprechenden Ereignisse gefunden. Dies führt im Raum der Oszillationsparameter zu verbotenen Bereichen, welche sich in Abbildung 3.3 rechts von den verschiedenen Linien befinden (90 % CL). Insbesondere das KARMEN-Experiment und sein Nachfolger KARMEN 2 am Rutherford-Appleton-Laboratorium, welche im wesentlichen den gleichen Aufbau wie LSND haben (mit $`L=17.6`$ m), schließen einen großen Teil des LSND-Bereichs aus. Der Grund für diese widersprüchlichen Resultate ist bis heute nicht bekannt. Allerdings wird das für 2001 am Fermilab in Planung befindliche Szintillator-Experiment MiniBooNE aufgrund einer deutlich verbesserten Statistik den gesamten von LSND erlaubten Parameterbereich untersuchen und somit die Frage nach $`\overline{\nu }_\mu \overline{\nu }_e`$-Oszillationen abschließend beantworten können.
#### 3.2.4 Analyse für drei Neutrino-Arten
In Tabelle 3.3 sind die Wertebereiche von $`\mathrm{\Delta }m^2`$ zusammengefaßt, welche eine Oszillationslösung der jeweiligen Neutrino-Anomalie ermöglichen. Da sich die zulässigen Wertebereiche je zweier $`\mathrm{\Delta }m^2`$ in keinem Fall überlappen und $`n`$ verschiedene Neutrinomassen $`n1`$ unabhängige $`\mathrm{\Delta }m^2`$ liefern, können nur dann alle drei Anomalien durch Neutrino-Oszillationen erklärt werden, wenn es vier leichte Neutrinos gibt.
Es sind verschiedene Modelle mit einem vierten, im Rahmen des SM nicht wechselwirkenden, Neutrino $`\nu _s`$ untersucht worden (Referenzen in ); da aber die Resultate von LSND und KARMEN bisher widersprüchlich sind, wird in der vorliegenden Arbeit die LSND-Anomalie nicht weiter berücksichtigt und von drei leichten Neutrinos ausgegangen.
Analysiert man alle Daten aus den Oszillationslösungen der übrigen beiden Neutrinodefizite und aus den Experimenten, welche keine Hinweise auf Oszillationen liefern, unter der Annahme dreier miteinander mischender leichter Neutrinos, so erhält man je nach Erklärung der Sonnen-Neutrino-Anomalie folgende erlaubten Bereiche für die Massen und Mischungen :
MSW-Effekt (kleiner Winkel):
$`0.410^5\text{eV}^2\mathrm{\Delta }m_{\text{sun}}^2\mathrm{\hspace{0.33em}1.2}10^5\text{eV}^2`$ (3.20)
$`0.410^3\text{eV}^2\mathrm{\Delta }m_{\text{atm}}^2\mathrm{\hspace{0.33em}8.0}10^3\text{eV}^2`$ (3.21)
$`33\mathrm{\Delta }m_{\text{atm}}^2/\mathrm{\Delta }m_{\text{sun}}^2\mathrm{\hspace{0.33em}2000}`$ (3.22)
$`|𝐔|=\left(\begin{array}{ccc}1& 0.030.05& 1\\ 0.020.05& 0.710.87& 0.490.71\\ 0.010.04& 0.480.71& 0.710.87\end{array}\right)`$ (3.26)
MSW-Effekt (großer Winkel):
$`0.810^5\text{eV}^2\mathrm{\Delta }m_{\text{sun}}^2\mathrm{\hspace{0.33em}3.0}10^5\text{eV}^2`$ (3.27)
$`0.410^3\text{eV}^2\mathrm{\Delta }m_{\text{atm}}^2\mathrm{\hspace{0.33em}8.0}10^3\text{eV}^2`$ (3.28)
$`13\mathrm{\Delta }m_{\text{atm}}^2/\mathrm{\Delta }m_{\text{sun}}^2\mathrm{\hspace{0.33em}1000}`$ (3.29)
$`|𝐔|=\left(\begin{array}{ccc}0.870.94& 0.350.49& 1\\ 0.250.43& 0.610.82& 0.490.71\\ 0.170.35& 0.420.66& 0.710.87\end{array}\right)`$ (3.33)
Vakuumoszillationen:
$`0.610^{10}\text{eV}^2\mathrm{\Delta }m_{\text{sun}}^2\mathrm{\hspace{0.33em}1.1}10^{10}\text{eV}^2`$ (3.34)
$`0.410^3\text{eV}^2\mathrm{\Delta }m_{\text{atm}}^2\mathrm{\hspace{0.33em}6.0}10^3\text{eV}^2`$ (3.35)
$`310^6\mathrm{\Delta }m_{\text{atm}}^2/\mathrm{\Delta }m_{\text{sun}}^2\mathrm{\hspace{0.33em}10}^8`$ (3.36)
$`|𝐔|=\left(\begin{array}{ccc}0.710.88& 0.480.71& 1\\ 0.340.61& 0.500.76& 0.510.71\\ 0.240.50& 0.360.62& 0.710.86\end{array}\right)`$ (3.40)
Ein realistisches Massenmodell für die Fermionen muß demnach auf Neutrino-Eigenschaften führen, welche mit einem der drei obigen Fälle innerhalb der Grenzen konsistent sind.
## Kapitel 4 Das $`SO(10)`$-Massenmodell
### 4.1 Bestimmung der Symmetriebrechungsskalen
Gegenstand dieser Arbeit ist die detaillierte Analyse eines Modellansatzes für die fermionischen Massenmatrizen im Rahmen einer nichtsupersymmetrischen $`SO(10)`$-GUT mit dem Brechungsschema
$$SO(10)\stackrel{M_U}{}G_I\stackrel{M_I}{}G_{\text{SM}}\stackrel{M_Z}{}SU(3)_CU(1)_{\text{em}}$$
(4.1)
Als intermediäre Symmetriegruppe $`G_I`$ sollen zunächst sowohl $`G_{\text{PS}}`$ als auch $`G_{\text{PS}}D`$ in Frage kommen. Deren Vorteil gegenüber Symmetriegruppen, welche $`SU(3)_CU(1)_{BL}`$ enthalten, besteht in der Vereinheitlichung von Quarks und Leptonen bereits bei der Skala $`M_I`$. Das hat unter anderem zur Folge, daß die $`SO(10)`$-Relationen (2.35-2.40) zwischen den verschiedenen Massenmatrizen nicht nur oberhalb von $`M_U`$, sondern auch im Energiebereich zwischen $`M_I`$ und $`M_U`$ gelten. Ferner resultiert aus dem links-rechts-symmetrischen Fermionspektrum (2.3.2) auf natürliche Weise die Existenz rechtshändiger Neutrinos.
Gemäß Tabelle 2.2 wird die Symmetriebrechung bei $`M_U`$ durch eine 210 beziehungsweise 54 realisiert; die Brechung bei $`M_I`$ erfolgt in beiden Fällen durch das SM-Singulett einer $`(\mathbf{10,1,3})_{126}`$. Bezüglich der Massen der Higgs-Teilchen wird von der Gültigkeit der „Extended Survival Hypothesis“ ausgegangen, das heißt nur diejenigen Komponenten von $`SO(10)`$-Higgs-Darstellungen besitzen Massen $`mM_U`$, welche für die Symmetriebrechungen bei $`M_I`$ und $`M_Z`$ sowie für die Erzeugung der Fermionmassen erforderlich sind. Higgs-Teilchen mit Massen der Größenordnung $`M_I`$ können demnach in folgenden Darstellungen liegen:
* in der $`(\mathbf{10,1,3})_{126}`$, da sie die Symmetriebrechung bei $`M_I`$ realisiert.
* in $`(\mathbf{1,2,2})_{10/120}`$ und $`(\mathbf{15,2,2})_{120/126}`$, da diese für die Massenerzeugung der Fermionen in Frage kommen.
* in Modellen mit $`G_{\text{PS}}D`$-Symmetrie in der $`(\overline{\mathrm{𝟏𝟎}}\mathbf{,3,1})_{126}`$, welche aufgrund der $`D`$-Parität das Gegenstück zur $`(\mathbf{10,1,3})_{126}`$ bildet (siehe Abschnitt 2.3.2).
Wie bereits in Abschnitt 2.3.3 erläutert wurde, ist das SM-Higgs-Doublett eine Linearkombination der $`SU(2)_L`$-Doubletts in den $`(\mathbf{1,2,2})`$\- und $`(\mathbf{15,2,2})`$-Darstellungen.
Um ein realistisches $`SO(10)`$-Massenmodell konstruieren zu können, ist wegen (2.41) ein Higgs-Spektrum erforderlich, welches über den einfachsten Fall einer komplexen 10 hinausgeht. Während im Bereich unterhalb von $`M_I`$ der Teilcheninhalt stets der des SM mit einem Higgs-Doublett sein soll, ist die Anzahl der für die Erzeugung der Fermionmassen relevanten $`G_{\text{PS}}`$-Higgs-Darstellungen $`(\mathbf{1,2,2})_{10/120}`$ und $`(\mathbf{15,2,2})_{120/126}`$ mit Massen $`M_I`$ zunächst nicht weiter festgelegt und kann an die konkreten Anforderungen des Modells angepaßt werden. Der Umfang des Teilchenspektrums bei $`M_I`$ hat allerdings direkte Auswirkungen auf die Skalenabhängigkeit der $`G_{\text{PS}}[D]`$-Eichkopplungen, wie man an den Renormierungsgruppengleichungen (B.21-B.23) erkennt. Bedeutsam wird dieser Einfluß, wenn man die phänomenologischen Eigenschaften des Massenmodells wie zum Beispiel die Zerfallsraten der Nukleonen berechnen will. Dazu müssen nämlich die Werte der Symmetriebrechungsskalen $`M_{U,I}`$ und die der verschiedenen Eichkopplungen bei diesen Skalen bekannt sein, und um diese quantitativ bestimmen zu können, ist eine numerische Integration der Renormierungsgruppengleichungen für die Kopplungen sowohl des SM (B.7-B.9) als auch des $`G_{\text{PS}}[D]`$-Modells (B.21-B.23) erforderlich.
Im folgenden werden zunächst $`M_{U,I}`$ und die Kopplungswerte für beide intermediären Symmetriegruppen in Abhängigkeit vom Higgs-Spektrum numerisch bestimmt, wobei das verwendete Verfahren dem in entspricht. Dort ist der minimale Fall mit einer 10 untersucht worden; in wurde der Einfluß des Higgs-Spektrums auf die $`\beta `$-Funktionen der $`G_I`$-Kopplungen analysiert, ohne jedoch die Renormierungsgruppengleichungen numerisch zu lösen.
#### 4.1.1 $`SU(4)_CSU(2)_LSU(2)_R`$-Modell
Die Parameter $`\mathrm{\Delta }_{R,L}`$ in (B.21-B.23) haben hier die festen Werte $`\mathrm{\Delta }_R=1`$ und $`\mathrm{\Delta }_L=0`$, die Werte von $`N_1`$ und $`N_{15}`$ können frei vorgegeben werden. Unter Berücksichtigung des Prinzips, daß die Anzahl der Higgs-Teilchen bei gegebenem Massenmodell grundsätzlich so klein wie möglich gewählt werden soll, erscheint eine Untersuchung der Fälle $`N_1=1,\mathrm{}\mathrm{,5}`$ und $`N_{15}=\mathrm{1,2,3}`$ sinnvoll. Die Vorgehensweise sieht dann wie folgt aus:
* Für $`M_I`$ wird ein Wert vorgegeben, der sinnvollerweise zwischen $`M_Z`$ und $`M_{\text{Planck}}10^{19}`$ GeV liegen kann.
* Ausgehend von den bekannten Werten der SM-Eichkopplungen bei $`M_Z`$ (siehe Tabelle C.1) werden die Renormierungsgruppengleichungen (B.7-B.9) unter Vernachlässigung der Yukawabeiträge von $`M_Z`$ bis $`M_I`$ integriert.
* Aus den nun bekannten SM-Kopplungen bei $`M_I`$ werden die Kopplungen des $`G_{\text{PS}}`$-Modells bei dieser Skala berechnet. Die Anschlußbedingungen lauten
$`\alpha _{4C}^1(M_I)`$ $`=`$ $`\alpha _3^1(M_I)+{\displaystyle \frac{1}{12\pi }}`$ (4.2)
$`\alpha _{2L}^1(M_I)`$ $`=`$ $`\alpha _2^1(M_I)`$ (4.3)
$`\alpha _{2R}^1(M_I)`$ $`=`$ $`{\displaystyle \frac{5}{3}}\alpha _1^1(M_I){\displaystyle \frac{2}{3}}\alpha _3^1(M_I)+{\displaystyle \frac{1}{3\pi }}`$ (4.4)
Die Faktoren 5/3 und 2/3 in (4.4) stammen aus der korrekten Normierung der $`G_{\text{PS}}`$-Generatoren im Rahmen der übergeordneten $`SO(10)`$ , während die Korrekturen $`1/(12\pi )`$ und $`1/(3\pi )`$ auf Schwelleneffekte zurückzuführen sind . Diese Effekte modifizieren die „naiven“ Anschlußbedingungen in der zweiten Ordnung der Störungsrechnung. Vernachlässigt man logarithmische Korrekturen $`\mathrm{ln}(m_i/M_{I,U})`$, wobei $`m_i`$ für die in der Regel nicht exakt bekannten Massen aller Teilchen steht, welche ihre Massen durch die Symmetriebrechung bei $`M_{I,U}`$ erhalten, können die Anschlußbedingungen allgemein als
$$\alpha _j^1(M_{I/U})\frac{1}{12\pi }S_2(𝒢_j)=\alpha _k^1(M_{I/U})\frac{1}{12\pi }S_2(𝒢_k)$$
(4.5)
geschrieben werden . $`S_2(𝒢_j)`$ bezeichnet hier den Dynkin-Index der adjungierten Darstellung der Eichgruppe $`G_j`$ (siehe Anhang A). In Abschnitt 4.2 werden die Ursachen von Schwellenkorrekturen und ihr Einfluß auf die Bestimmung der Symmetriebrechungsskalen ausführlich diskutiert.
* Ausgehend von den $`G_{\text{PS}}`$-Kopplungen bei $`M_I`$ integriert man die Renormierungsgruppengleichungen (B.21-B.23) von $`M_I`$ zu höheren Energien und überprüft, ob eine Skala $`\mu M_U>M_I`$ existiert, bei welcher sich die drei Kopplungen unter Berücksichtigung der ebenfalls durch Schwelleneffekte modifizierten Anschlußbedingungen
$`\alpha _U^1(M_U)`$ $`=`$ $`\alpha _{4C}^1(M_U)+{\displaystyle \frac{1}{3\pi }}`$ (4.6)
$`=`$ $`\alpha _{2L}^1(M_U)+{\displaystyle \frac{1}{2\pi }}`$ (4.7)
$`=`$ $`\alpha _{2R}^1(M_U)+{\displaystyle \frac{1}{2\pi }}`$ (4.8)
in einem Punkt treffen. Ist das der Fall, so sind die Brechungsskalen $`M_I`$ und $`M_U`$ sowie die Werte der verschiedenen Kopplungskonstanten bei $`M_{I,U}`$ bekannt.
Auf diese Weise ist für jedes der in Frage kommenden Higgs-Spektren der gesamte erlaubte Parameterbereich für $`M_I`$ zu untersuchen; dabei besteht durchaus die Möglichkeit, daß keine Vereinheitlichung erreicht werden kann.
Die Ergebnisse sind in Anhang C.2 zusammengefaßt. Man erkennt zunächst, daß für $`N_{15}=3`$ und $`N_1=1,\mathrm{}\mathrm{,5}`$ sowie $`N_{15}=2`$ und $`N_1=\mathrm{1,2,3}`$ keine Vereinheitlichung stattfindet. In den übrigen untersuchten Fällen sind folgende Eigenschaften erkennbar:
* Bei festem $`N_{15}`$ sind $`M_U`$ und $`\alpha _U`$ umso kleiner und $`M_I`$ umso größer, je größer $`N_1`$ ist.
* Bei festem $`N_1`$ sind $`M_U`$ und $`\alpha _U`$ umso größer und $`M_I`$ umso kleiner, je größer $`N_{15}`$ ist.
Eine qualitative Abschätzung der Protonlebensdauer gemäß (2.9) führt unter Berücksichtigung der experimentellen Grenzen mit $`\alpha _U1/35\mathrm{}1/20`$ zu der Einschränkung $`M_U310^{15}`$ GeV für die Vereinheitlichungsskala. Deswegen sind die Modelle mit $`(N_1,N_{15})=(\mathrm{4,1})`$ und $`(\mathrm{5,1})`$ wegen ihrer relativ kleinen Werte von $`M_U`$ im Vergleich zu den übrigen vom phänomenologischen Standpunkt her eher ungeeignet.
#### 4.1.2 $`SU(4)_CSU(2)_LSU(2)_RD`$-Modell
Die Parameter $`\mathrm{\Delta }_{R,L}`$ haben in diesem Fall wegen des zwangsläufig links-rechts-symmetrischen Teilcheninhalts beide den Wert 1. Im Vergleich zum $`G_{\text{PS}}`$-Modell ergibt sich hier ein wesentlicher Unterschied in der Vorgehensweise, der auf den Anschlußbedingungen bei $`M_I`$ beruht. Zusätzlich zu (4.2-4.4) gilt in Modellen mit $`D`$-Parität nämlich $`\alpha _{2L}(\mu )=\alpha _{2R}(\mu )`$ für alle $`\mu `$ mit $`M_I\mu M_U`$. Das führt auf
$`\alpha _{4C}^1(M_I)`$ $`=`$ $`\alpha _3^1(M_I)+{\displaystyle \frac{1}{12\pi }}`$ (4.9)
$`\alpha _{2L}^1(M_I)`$ $`=`$ $`\alpha _2^1(M_I)=\alpha _{2R}^1(M_I)`$ (4.10)
$`0`$ $`=`$ $`\alpha _2^1(M_I){\displaystyle \frac{5}{3}}\alpha _1^1(M_I)+{\displaystyle \frac{2}{3}}\alpha _3^1(M_I){\displaystyle \frac{1}{3\pi }}`$ (4.11)
Die letzte Gleichung liefert eine Einschränkung für die SM-Kopplungen, welche bei genau einer Skala $`M_I`$ erfüllt ist. In ist gezeigt worden, daß dieser Wert für $`M_I`$ von den physikalischen Eigenschaften des Modells bei Skalen $`\mu >M_I`$ völlig unabhängig ist; insbesondere hat das Higgs-Spektrum keinen Einfluß auf $`M_I`$ und die Kopplungskonstanten bei $`\mu =M_I`$. In Tabelle C.2 sind die Werte dieser Größen angegeben. Es fällt auf, daß der Wert von $`M_I`$ um etwa zwei bis drei Größenordnungen über den entsprechenden Werten im $`G_{\text{PS}}`$-Modell liegt.
Das Problem reduziert sich also auf die Bestimmung von $`M_U`$ und $`\alpha _U`$ in Abhängigkeit vom Teilcheninhalt. Von den zwei bekannten Kopplungen bei $`M_I`$ ausgehend werden (B.21-B.22) bis zu der Skala $`\mu M_U`$ integriert, bei welcher $`\alpha _{4C}`$ und $`\alpha _{2L}`$ die Bedingungen
$`\alpha _U^1(M_U)`$ $`=`$ $`\alpha _{4C}^1(M_U)+{\displaystyle \frac{1}{3\pi }}`$ (4.12)
$`=`$ $`\alpha _{2L}^1(M_U)+{\displaystyle \frac{1}{2\pi }}`$ (4.13)
erfüllen. Die Resultate sind in Tabelle C.3 zusammengefaßt, es sind folgende Tendenzen in den Lösungen erkennbar:
* Bei festem $`N_{15}`$ sind $`M_U`$ und $`\alpha _U`$ umso kleiner, je größer $`N_1`$ ist.
* Bei festem $`N_1`$ sind $`M_U`$ und $`\alpha _U`$ umso größer, je größer $`N_{15}`$ ist.
Für alle untersuchten Fälle findet eine Vereinheitlichung der Kopplungen statt, allerdings bei vergleichsweise kleinen $`M_U`$-Werten von $`(12)10^{15}`$ GeV. Das führt zu Vorhersagen für die Lebensdauern der Nukleonen, die den experimentellen Grenzen widersprechen, weshalb $`G_{\text{PS}}D`$ als intermediäre Symmetriegruppe ausscheidet. Im folgenden ist $`G_I=G_{\text{PS}}`$; der Higgs-Inhalt des Modells wird später aus den Anforderungen des Ansatzes für die Massenmatrizen unter Berücksichtigung der Resultate in Abschnitt 4.1.1 bestimmt.
### 4.2 Schwellenkorrekturen
In diesem Abschnitt wird auf den Ursprung der Schwellenkorrekturen und deren Effekte auf die Bestimmung von Symmetriebrechungsskalen eingegangen.
Wird eine auf der Gruppe $`G_J`$ beruhende Eichsymmetrie bei einer Skala $`M`$ spontan in die zur Untergruppe $`G_j`$ gehörige Symmetrie gebrochen, ist damit formal der Übergang von der vollen Eichtheorie zu einer bei Energien $`\mu <M`$ gültigen effektiven Theorie verbunden . Diese effektive Niederenergie-Näherung erhält man durch Ausintegration der schweren Freiheitsgrade aus dem Wirkungsfunktional. Wenn $`\mathrm{\Phi }`$ die Felder bezeichnet, die zu Teilchen mit Massen der Größenordnung $`M`$ gehören, und $`\varphi `$ für die auch unterhalb von $`M`$ masselosen Felder steht, erhält man die Wirkung $`S_j[\varphi ]`$ der effektiven Theorie aus der ursprünglichen Wirkung $`S_J[\varphi ,\mathrm{\Phi }]`$ durch Funktionalintegration über die schweren Felder:
$$\mathrm{exp}(iS_j[\varphi ])=[d\mathrm{\Phi }]\mathrm{exp}(iS_J[\varphi ,\mathrm{\Phi }])$$
(4.14)
Hierbei ist die Invarianz von $`S_j[\varphi ]`$ unter $`G_j`$-Transformationen zu fordern. Um die Integration in (4.14) ausführen zu können, muß in der Wirkung $`S_J[\varphi ,\mathrm{\Phi }]`$ ein Eichfixierungsterm $`f_J(\varphi ,\mathrm{\Phi })`$ eingeführt werden. Nach der Ausintegration der schweren Felder verbleibt von $`f_J(\varphi ,\mathrm{\Phi })`$ jedoch ein Rest, der die $`G_j`$-Eichinvarianz der neuen Wirkung $`S_j[\varphi ]`$ verletzt. Dieses Problem kann durch eine Redefinition der Kopplungskonstanten und der Eichfelder in $`S_j`$ behoben werden. Als Konsequenz ändert sich die naive Anschlußbedingung $`\alpha _j(M)=\alpha _J(M)`$ in
$$\alpha _j^1(M)=\alpha _J^1(M)\lambda _j(M),$$
(4.15)
wobei $`\lambda _j(M)`$ in erster Ordnung der Störungsrechung durch
$`\lambda _j(M)`$ $`=`$ $`{\displaystyle \frac{1}{12\pi }}(S_2(_j^{\left(G\right)})21S_2(_j^{\left(G\right)})\mathrm{ln}(m_G/M)`$ (4.16)
$`+\mathrm{\Lambda }{\displaystyle \underset{S}{}}S_2(_j^{\left(S\right)})\mathrm{ln}(m_S/M)+8{\displaystyle \underset{F}{}}S_2(_j^{\left(F\right)})\mathrm{ln}(m_F/M))`$
gegeben ist . Hierbei ist $`_j^{\left(G\right)}`$ die Darstellung von $`G_j`$, nach welcher die im Rahmen der Symmetriebrechung massiv gewordenen Eichbosonen von $`G_J`$ transformieren. Weiterhin steht $`_j^{(F,S)}`$ für die $`𝒢_j`$-Darstellungen, in denen die Fermionen beziehungsweise Higgs-Teilchen mit Massen $`m_{F,S}M`$ liegen; $`S_2`$ ist der Dynkin-Index dieser Darstellungen (siehe Tabelle A.1). Der Operator $`\mathrm{\Lambda }`$ projeziert die Goldstone-Bosonen aus dem Higgs-Spektrum heraus. $`\lambda _j(M)`$ hat seinen physikalischen Ursprung in den Strahlungskorrekturen zum Propagator der Eichbosonen der effektiven Theorie durch die schweren Teilchen.
In der Analyse des letzten Abschnitts wurde bei den Anschlußbedingungen lediglich der erste, nicht von den Massen abhängige, Term von $`\lambda _j`$ in abgewandelter Form benutzt. Das entspricht der üblicherweise gemachten Annahme, daß die durch Symmetriebrechung bei einer Skala $`M`$ erzeugten Teilchenmassen sich nur minimal von $`M`$ unterscheiden und die logarithmischen Terme vernachlässigt werden können.
Hier soll nun der Einfluß der zu den Skalaren gehörigen logarithmischen Terme auf die Bestimmung der Werte der Symmetriebrechungsskalen $`M_I`$ und $`M_U`$ untersucht werden, wenn man für die Argumente der Logarithmen einen Wertebereich von 0.1 bis 10 zuläßt. Das ist angesichts der Komplexität des Higgs-Potentials in Modellen mit nichtminimalem Higgs-Inhalt und der damit verbundenen großen Anzahl unbekannter Koeffizienten der Größenordnung 1 durchaus plausibel. Das Verfahren lehnt sich an an, wo dies für $`SO(10)`$-Modelle mit minimalem Higgs-Inhalt durchgeführt wurde. Beiträge durch Fermionen gibt es bei $`M_{U,I}`$ nicht, da die rechtshändigen Neutrinos als einzige schwere Teilchen SM-Singuletts sind. Die Eichbosonbeiträge werden vernachlässigt, da der Higgs-Inhalt der Theorie deutlich umfangreicher als der Eichboson-Inhalt ist.
Die für die Berechnung der $`\lambda _j`$-Koeffizienten erforderlichen Dynkin-Indizes der einzelnen Higgs-Darstellungen sind in Tabelle A.1 angegeben. Mit ihrer Hilfe lassen sich die $`\lambda _j^{I,U}`$ für das hier diskutierte Modell bestimmen; Anhang A.4 enthält alle relevanten Resultate. Der Einfluß der Schwellenkorrekturen auf die Vorhersagen von $`M_{U,I}`$ wird durch folgende Beziehungen beschrieben :
$`\mathrm{\Delta }\mathrm{ln}(M_U/M_Z)`$ $`=`$ $`{\displaystyle \frac{K_\lambda A_IJ_\lambda B_I}{A_UB_IA_IB_U}}`$ (4.17)
$`\mathrm{\Delta }\mathrm{ln}(M_I/M_Z)`$ $`=`$ $`{\displaystyle \frac{J_\lambda B_UK_\lambda A_U}{A_UB_IA_IB_U}}`$ (4.18)
Die Größen $`A_{I,U}`$ und $`B_{I,U}`$ sind Linearkombinationen der führenden Koeffizienten der Eichkopplungs-$`\beta `$-Funktionen
$`A_U`$ $`=`$ $`2\beta _{4C}\beta _{2L}\beta _{2R}`$ (4.19)
$`B_U`$ $`=`$ $`{\displaystyle \frac{2}{3}}\beta _{4C}+{\displaystyle \frac{5}{3}}\beta _{2L}\beta _{2R}`$ (4.20)
$`A_I`$ $`=`$ $`{\displaystyle \frac{8}{3}}\beta _3\beta _2{\displaystyle \frac{5}{3}}\beta _1A_U`$ (4.21)
$`B_I`$ $`=`$ $`{\displaystyle \frac{5}{3}}(\beta _2\beta _1)B_U`$ (4.22)
und können Anhang B entnommen werden, $`J_\lambda `$ und $`K_\lambda `$ sind als
$`J_\lambda `$ $`=`$ $`2\pi \left(2\lambda _{4C}^U+\lambda _{2L}^U+\lambda _{2R}^U{\displaystyle \frac{8}{3}}\lambda _{3C}^I+\lambda _{2L}^I+{\displaystyle \frac{5}{3}}\lambda _{1Y}^I\right)`$ (4.23)
$`K_\lambda `$ $`=`$ $`2\pi \left({\displaystyle \frac{2}{3}}\lambda _{4C}^U{\displaystyle \frac{5}{3}}\lambda _{2L}^U+\lambda _{2R}^U{\displaystyle \frac{5}{3}}\lambda _{2L}^I+{\displaystyle \frac{5}{3}}\lambda _{1Y}^I\right)`$ (4.24)
definiert. Dabei ist $`\lambda _j^{I,U}`$ als Summe der Beiträge der einzelnen Higgs-Darstellungen zu verstehen.
Um daraus quantitative Resultate zu erhalten, sei an dieser Stelle das Higgs-Spektrum des noch zu konstruierenden Massenmodells vorweggenommen. Es wird sich abgesehen von $`\mathrm{\Delta }_R=1`$ und $`\mathrm{\Delta }_L=0`$ zu $`(N_1,N_{15})=(\mathrm{4,2})`$ ergeben, woraus sich für die durch Schwelleneffekte verursachten Unsicherheiten in den Vorhersagen für $`M_{U,I}`$
$`\mathrm{\Delta }\mathrm{ln}(M_U/M_Z)`$ $`=`$ $`{\displaystyle \frac{42}{215}}\eta _{\left(\mathbf{1,2,2}\right)}+{\displaystyle \frac{24}{215}}\eta _{\left(\mathbf{15,2,2}\right)}+{\displaystyle \frac{141}{430}}\eta _{\left(\mathbf{10,1,3}\right)}`$ (4.25)
$`+{\displaystyle \frac{12}{215}}\eta _{\mathrm{𝟐𝟏𝟎}}{\displaystyle \frac{17}{43}}\eta _{\mathrm{𝟏𝟐𝟔}}+{\displaystyle \frac{12}{215}}\eta _{\mathrm{𝟏𝟐𝟎}}+{\displaystyle \frac{24}{215}}\eta _{\mathrm{𝟏𝟎}}`$
$`\mathrm{\Delta }\mathrm{ln}(M_I/M_Z)`$ $`=`$ $`{\displaystyle \frac{7}{43}}\eta _{\left(\mathbf{1,2,2}\right)}{\displaystyle \frac{4}{43}}\eta _{\left(\mathbf{15,2,2}\right)}{\displaystyle \frac{44}{43}}\eta _{\left(\mathbf{10,1,3}\right)}`$ (4.26)
$`{\displaystyle \frac{2}{43}}\eta _{\mathrm{𝟐𝟏𝟎}}+{\displaystyle \frac{50}{43}}\eta _{\mathrm{𝟏𝟐𝟔}}{\displaystyle \frac{2}{43}}\eta _{\mathrm{𝟏𝟐𝟎}}{\displaystyle \frac{4}{43}}\eta _{\mathrm{𝟏𝟎}}`$
ergibt. $`\eta _{}`$ steht für $`\mathrm{ln}(m_{}/M_U)`$ ($`=\mathbf{10,120,126,210}`$) beziehungsweise $`\mathrm{ln}(m_{}/M_I)`$ ($`=(\mathbf{1,2,2}),(\mathbf{15,2,2}),(\mathbf{10,1,3})`$), wobei vereinfachend davon ausgegangen wird, daß alle Teilchen in einer Darstellung $``$, welche ihre Massen durch den Symmetriebrechungsschritt bei $`M_{U,I}`$ erhalten, dieselbe Masse haben. Läßt man nun, wie oben angekündigt, für die $`\eta _{}`$ Werte zwischen $`\mathrm{ln}(0.1)`$ und $`\mathrm{ln}(10)`$ zu, erhält man folgende Maximalbeträge für $`\mathrm{\Delta }\mathrm{ln}(M_{U,I}/M_Z)`$:
$`|\mathrm{\Delta }\mathrm{ln}(M_U/M_Z)|_{\text{max}}`$ $`=`$ $`{\displaystyle \frac{539}{430}}\mathrm{ln}(10)=\mathrm{\hspace{0.33em}2.89}`$ (4.27)
$`|\mathrm{\Delta }\mathrm{ln}(M_I/M_Z)|_{\text{max}}`$ $`=`$ $`{\displaystyle \frac{113}{43}}\mathrm{ln}(10)=\mathrm{\hspace{0.33em}6.05}`$ (4.28)
Demnach können die Schwellenkorrekturen im vorliegenden Fall aufgrund ihres Einflusses auf die Anschlußbedingungen gemäß (4.15) den Wert von $`M_U`$ um einen Faktor $`10^{\pm 1.25}`$ und den von $`M_I`$ um einen Faktor $`10^{\pm 2.63}`$ modifizieren.
Obwohl hier lediglich eine Abschätzung des Maximaleffekts durchgeführt wurde, sollte bei den Angaben über die Werte von Symmetriebrechungsskalen der mögliche Einfluß der unbekannten Higgs-Massen nicht vergessen werden.
### 4.3 Der Ansatz für die Massenmatrizen
Im Rahmen des SM ist die Form der fermionischen Massenmatrizen in keiner Weise eingeschränkt. Es besteht aber allgemein die Überzeugung, daß eine fundamentalere Theorie als das SM existiert, welche auch den Fermionsektor einschließlich der Struktur der Massenmatrizen erklären kann. In Ermangelung einer solchen Theorie ist bis heute eine Vielzahl von phänomenologisch motivierten Ansätzen für die Massenmatrizen vorgeschlagen worden. Sie zeichnen sich durch bestimmte Symmetrieeigenschaften wie zum Beispiel Hermitezität und sogenannte Texturen, das heißt Nullen als Matrixeinträge an bestimmten Stellen, aus. Der Grundgedanke besteht darin, die Anzahl der freien Parameter in den Matrizen kleiner als die Zahl der zu reproduzierenden observablen Massen und Mischungen zu machen, um Vorhersagen über Beziehungen zwischen diesen Größen zu erhalten und somit möglicherweise etwas über die Eigenschaften der zugrundeliegenden Theorie zu erfahren.
Der wohl populärste dieser Ansätze geht auf Fritzsch zurück und basiert auf zwei Annahmen: Die Quark-Massenmatrizen sind hermitesch und haben die „Nearest Neighbour Interaction“ (NNI)-Form. Letzteres bedeutet, daß nur das schwerste Fermion einer Art seine Masse direkt über die Yukawakopplung erhält; die Massen der leichten Teilchen kommen durch Wechselwirkungen respektive Mischungen mit ihren nächsten Nachbarn in der Massenmatrix zustande. Für zwei und drei Generationen haben die Matrizen demnach folgende Gestalt:
$$𝐌=\left(\begin{array}{ccc}0\hfill & A\hfill & \\ A^{}\hfill & B\hfill & \end{array}\right)\text{beziehungsweise}𝐌=\left(\begin{array}{ccc}0\hfill & A\hfill & 0\hfill \\ A^{}\hfill & 0\hfill & B\hfill \\ 0\hfill & B^{}\hfill & C\hfill \end{array}\right)$$
(4.29)
Wendet man den Fritzsch-Ansatz auf die Quark-Massenmatrizen im Falle zweier Generationen an, erhält man für den Cabibbo-Winkel die mit $`\varphi \pi /2`$ sehr gute Vorhersage
$$\theta _C|\sqrt{m_d/m_s}e^{i\varphi }\sqrt{m_u/m_c}|$$
(4.30)
Ferner läßt sich mit $`|A||B|`$ auch die Massenhierarchie erklären.
Problematisch ist der Ansatz jedoch für drei Fermiongenerationen, da es keine Lösungen gibt, welche gleichzeitig die große $`t`$-Quarkmasse und den kleinen Wert für das CKM-Matrixelement $`V_{cb}`$ realisieren können . Daran ändert sich auch nichts, wenn man den Fritzsch-Ansatz in eine $`SO(10)`$-GUT ohne oder mit Supersymmetrie einbettet.
Eine Möglichkeit, dieses Problem zu umgehen, könnte darin bestehen, hermitesche Massenmatrizen mit einer anderen Struktur als (4.29) zu verwenden. In ist eine systematische Analyse aller Kombinationen von hermiteschen Quark-Massenmatrizen mit insgesamt $`5`$ unabhängigen Null-Einträgen durchgeführt worden, aber auch dort sind keine wirklich überzeugenden Lösungen gefunden worden.
Es deutet also einiges darauf hin, daß ein phänomenologisch erfolgreiches Massenmodell nichthermitesche Matrizen enthalten muß. Hier hat sich das Interesse auf supersymmetrische GUTs konzentriert, wobei zur Massenerzeugung häufig von nichtrenormierbaren Operatoren Gebrauch gemacht wird, welche aus einer noch fundamentaleren Theorie stammen sollen; bietet einen Überblick über solche Ansätze in SUSY-GUTs.
Modelle mit nichthermiteschen Massenmatrizen, welche die NNI-Form besitzen, haben im Rahmen supersymmetrischer $`SU(5)`$\- und $`SO(10)`$-Theorien vergleichsweise gute Resultate auch im Neutrinosektor geliefert. Daß ein solcher Ansatz auch in einem $`SO(10)`$-Modell ohne Supersymmetrie erfolgreich sein kann, wird in dieser Arbeit gezeigt werden. Ein früherer Versuch in dieser Richtung war aufgrund der zu einfachen Struktur des verwendeten Higgs-Spektrums (eine $`\mathrm{𝟏𝟎}`$ und eine $`\mathrm{𝟏𝟐𝟎}`$) gescheitert .
Ein weiteres Argument zugunsten nichthermitescher NNI-Matrizen liefert . Dort ist gezeigt worden, daß ein NNI-Ansatz für die Quark-Massenmatrizen im SM keine physikalischen Konsequenzen beinhaltet, sofern die Zahl der Fermiongenerationen nicht größer als vier ist. Im SM kann man beide Matrizen unabhängig von deren Ausgangsgestalt durch Transformationen auf NNI-Form bringen, ohne daß sich observable Größen ändern. Der Grund dafür liegt in der Tatsache, daß von den vier Mischungsmatrizen im Quarksektor $`𝐔_{L,R}`$ und $`𝐃_{L,R}`$ aus (1.29) nur die Kombination $`𝐕𝐔_L^{}𝐃_L`$ physikalisch relevant ist; insbesondere die rechtshändigen Mischungen sind experimentell nicht beobachtbar. Es handelt sich beim NNI-Ansatz im SM also streng genommen nicht um einen Ansatz, sondern um eine spezielle Wahl der Basis im Raum der schwachen Eigenzustände. Wird der NNI-Ansatz jedoch in Theorien jenseits des SM eingebettet, ergeben sich daraus direkte Konsequenzen, die zumindest prinzipiell im Experiment überprüft werden können. So beeinflussen alle Mischungsmatrizen in (1.29) die Verzweigungsraten der Nukleonenzerfälle, und Beziehungen wie (2.35-2.40) liefern Eigenschaften des Neutrinosektors aus den Massenmatrizen der geladenen Fermionen.
Ausgangspunkt der nun folgenden Überlegungen wird die Annahme sein, daß die NNI-Form der Dirac-Massenmatrizen aus den Symmetrieeigenschaften einer wirklich grundlegenden Theorie der Elementarteilchen folgt. Diese nicht näher bekannte Theorie soll als Niederenergie-Näherung eine nichtsupersymmetrische $`SO(10)`$-GUT besitzen, welche wiederum über eine intermediäre $`G_{\text{PS}}`$-Symmetrie in das SM gebrochen wird. Für die fermionischen Massenmatrizen wird bei Skalen $`\mu M_I`$ der Ansatz
$$𝐌=\left(\begin{array}{ccc}0\hfill & A\hfill & 0\hfill \\ B\hfill & 0\hfill & C\hfill \\ 0\hfill & D\hfill & E\hfill \end{array}\right)$$
(4.31)
gemacht, wobei der Einfachheit halber nur reelle Einträge betrachtet werden. Das hat auf die Nukleonenzerfälle und Neutrinoeigenschaften, welche die zentralen Gegenstände der Untersuchung darstellen, keine nennenswerten Auswirkungen; lediglich auf das Problem der $`CP`$-Verletzung kann nicht eingegangen werden. Es sollte jedoch keine Schwierigkeiten bereiten, mit Hilfe komplexer Yukawa-Kopplungen die beobachtete $`CP`$-Verletzung zu reproduzieren.
Eine gängige Methode, bestimmte Texturen in den Massenmatrizen zu realisieren, besteht darin, eine globale $`U(1)`$-Familiensymmetrie zusätzlich zur Eichsymmetrie zu verwenden . Die drei $`SO(10)`$-Fermiondarstellungen transformieren sich dann gemäß $`\mathrm{𝟏𝟔}_j\mathrm{exp}(i\alpha _j\theta )\mathrm{𝟏𝟔}_j`$, wobei $`j`$ der Generationsindex ist und $`\alpha _j`$ die paarweise verschiedenen Ladungen bezüglich der $`U(1)`$ bezeichnet. Besitzt nun eine Higgs-Darstellung $`𝚽_k=\mathbf{10,120}`$ oder $`\mathrm{𝟏𝟐𝟔}`$ das $`U(1)`$-Transformationsverhalten $`𝚽_k\mathrm{exp}(i\beta _k\theta )𝚽_k`$, so sind nur solche Yukawa-Terme $`\mathrm{𝟏𝟔}_i\overline{𝚽}_k\mathrm{𝟏𝟔}_j`$ $`U(1)`$-invariant, welche die Bedingung $`\alpha _i+\alpha _j=\beta _k`$ erfüllen. Da die Einträge der Massenmatrizen gemäß
$$𝐌\left(\begin{array}{ccc}\alpha _1+\alpha _1\hfill & \alpha _1+\alpha _2\hfill & \alpha _1+\alpha _3\hfill \\ \alpha _1+\alpha _2\hfill & \alpha _2+\alpha _2\hfill & \alpha _2+\alpha _3\hfill \\ \alpha _1+\alpha _3\hfill & \alpha _2+\alpha _3\hfill & \alpha _3+\alpha _3\hfill \end{array}\right)$$
(4.32)
bestimmte $`U(1)`$-Ladungen besitzen, kann man die NNI-Form realisieren, indem nur Higgs-Darstellungen $`𝚽_k`$ mit den Ladungen $`\beta =\alpha _1+\alpha _2,\alpha _2+\alpha _3`$ und $`\alpha _3+\alpha _3`$ verwendet werden. Dabei ist es möglich, $`\alpha _1+\alpha _2=2\alpha _3`$ zu wählen, was die Anzahl der benötigten Darstellungen weiter reduziert. Die $`U(1)`$-Familiensymmetrie wird bei $`M_I`$ spontan gebrochen, da dort die massenerzeugenden Higgs-Darstellungen $`(\mathbf{1,2,2})`$ und $`(\mathbf{15,2,2})`$ selbst massiv werden. Unterhalb von $`M_I`$ werden aufgrund von Renormierungseffekten im allgemeinen auch die Einträge ungleich Null sein, welche gemäß (4.31) verschwinden.
Für die Yukawa-Kopplungsmatrizen kommen dann bei Skalen $`\mu M_I`$ unter Verwendung der Bezeichnungsweisen in Tabelle 2.3 folgende Möglichkeiten in Frage:
$`𝐘_{\mathrm{𝟏𝟎}}^{\left(1\right)}=\left(\begin{array}{ccc}0& x_1& 0\\ x_1& 0& 0\\ 0& 0& \stackrel{~}{x}_1\end{array}\right);𝐘_{\mathrm{𝟏𝟐𝟔}}^{\left(1\right)}=\left(\begin{array}{ccc}0& y_1& 0\\ y_1& 0& 0\\ 0& 0& \stackrel{~}{y}_1\end{array}\right);𝐘_{\mathrm{𝟏𝟐𝟎}}^{\left(1\right)}=\left(\begin{array}{ccc}0& z_1& 0\\ z_1& 0& 0\\ 0& 0& 0\end{array}\right);`$
$`𝐘_{\mathrm{𝟏𝟎}}^{\left(2\right)}=\left(\begin{array}{ccc}0& 0& 0\\ 0& 0& x_2\\ 0& x_2& 0\end{array}\right);𝐘_{\mathrm{𝟏𝟐𝟔}}^{\left(2\right)}=\left(\begin{array}{ccc}0& 0& 0\\ 0& 0& y_2\\ 0& y_2& 0\end{array}\right);𝐘_{\mathrm{𝟏𝟐𝟎}}^{\left(2\right)}=\left(\begin{array}{ccc}0& 0& 0\\ 0& 0& z_2\\ 0& z_2& 0\end{array}\right)`$ (4.33)
Das Higgs-Spektrum des Modells soll genau eine 126 enthalten, welche die Symmetriebrechung bei $`M_I`$ realisiert. Mehr als eine 126 zu verwenden ist nicht verboten, aber unnötig und würde die Vorhersagekraft insbesondere im Neutrinosektor beträchtlich reduzieren. Damit der See-Saw-Mechanismus gemäß (2.51) anwendbar ist, muß die Majorana-Massenmatrix der rechtshändigen Neutrinos (2.39) invertierbar sein. Aus diesem Grunde ist die Kopplung der 126 durch $`𝐘_{\mathrm{𝟏𝟐𝟔}}^{\left(1\right)}`$ aus (4.33) bestimmt.
Um die Asymmetrie gemäß (4.31) erzeugen zu können, sind beide 120-Darstellungen mit den Kopplungsmatrizen $`𝐘_{\mathrm{𝟏𝟐𝟎}}^{\left(1\right)}`$ und $`𝐘_{\mathrm{𝟏𝟐𝟎}}^{\left(2\right)}`$ erforderlich. Übernimmt man auch beide möglichen 10-Darstellungen aus (4.33) in den Higgs-Inhalt des Modells und beteiligt sämtliche $`(\mathbf{1,2,2})_{10/120}`$\- und $`(\mathbf{15,2,2})_{120/126}`$-Komponenten an der Massenerzeugung der Fermionen bei $`M_Z`$, das heißt gibt ihnen Massen $`M_I`$, so erhält man $`(N_1,N_{15})=(\mathrm{4,3})`$. Für diesen Fall findet gemäß der Analyse in Abschnitt 4.1.1 keine Vereinheitlichung statt, so daß eine der $`(\mathbf{15,2,2})`$ Massen der Größenordnung $`M_U`$ besitzen muß. Die Werte der Skalen und Kopplungen für den Fall $`N_1=4`$ und $`N_{15}=2`$ kann man Tabelle C.4 entnehmen; dabei erscheint $`M_U1.310^{16}`$ GeV bezüglich der Konsequenz für die Lebensdauer des Protons sinnvoll zu sein. Um entscheiden zu können, welche der $`(\mathbf{15,2,2})`$ für die Massenerzeugung die geringste Bedeutung hat, werden zunächst die Massenmatrizen unter Verwendung der vier $`(\mathbf{1,2,2})_{10/120}`$ und drei $`(\mathbf{15,2,2})_{120/126}`$ konstruiert. Deren allgemeine Gestalt sieht analog zu (2.35-2.40) bei Energien $`\mu M_I`$ dann folgendermaßen aus:
$`𝐌_d`$ $`=`$ $`\upsilon _d^{\left(1\right)}𝐘_{\mathrm{𝟏𝟎}}^{\left(1\right)}+\upsilon _d^{\left(2\right)}𝐘_{\mathrm{𝟏𝟎}}^{\left(2\right)}+\omega _d^{\left(1\right)}𝐘_{\mathrm{𝟏𝟐𝟔}}^{\left(1\right)}`$ (4.34)
$`+(\stackrel{~}{\upsilon }_d^{\left(1\right)}+\stackrel{~}{\omega }_d^{\left(1\right)})𝐘_{\mathrm{𝟏𝟐𝟎}}^{\left(1\right)}+(\stackrel{~}{\upsilon }_d^{\left(2\right)}+\stackrel{~}{\omega }_d^{\left(2\right)})𝐘_{\mathrm{𝟏𝟐𝟎}}^{\left(2\right)}`$
$`𝐌_e`$ $`=`$ $`\upsilon _d^{\left(1\right)}𝐘_{\mathrm{𝟏𝟎}}^{\left(1\right)}+\upsilon _d^{\left(2\right)}𝐘_{\mathrm{𝟏𝟎}}^{\left(2\right)}3\omega _d^{\left(1\right)}𝐘_{\mathrm{𝟏𝟐𝟔}}^{\left(1\right)}`$ (4.35)
$`+(\stackrel{~}{\upsilon }_d^{\left(1\right)}3\stackrel{~}{\omega }_d^{\left(1\right)})𝐘_{\mathrm{𝟏𝟐𝟎}}^{\left(1\right)}+(\stackrel{~}{\upsilon }_d^{\left(2\right)}3\stackrel{~}{\omega }_d^{\left(2\right)})𝐘_{\mathrm{𝟏𝟐𝟎}}^{\left(2\right)}`$
$`𝐌_u`$ $`=`$ $`\upsilon _u^{\left(1\right)}𝐘_{\mathrm{𝟏𝟎}}^{\left(1\right)}+\upsilon _u^{\left(2\right)}𝐘_{\mathrm{𝟏𝟎}}^{\left(2\right)}+\omega _u^{\left(1\right)}𝐘_{\mathrm{𝟏𝟐𝟔}}^{\left(1\right)}`$ (4.36)
$`+(\stackrel{~}{\upsilon }_u^{\left(1\right)}+\stackrel{~}{\omega }_u^{\left(1\right)})𝐘_{\mathrm{𝟏𝟐𝟎}}^{\left(1\right)}+(\stackrel{~}{\upsilon }_u^{\left(2\right)}+\stackrel{~}{\omega }_u^{\left(2\right)})𝐘_{\mathrm{𝟏𝟐𝟎}}^{\left(2\right)}`$
$`𝐌_\nu ^{\left(\text{Dir}\right)}`$ $`=`$ $`\upsilon _u^{\left(1\right)}𝐘_{\mathrm{𝟏𝟎}}^{\left(1\right)}+\upsilon _u^{\left(2\right)}𝐘_{\mathrm{𝟏𝟎}}^{\left(2\right)}3\omega _u^{\left(1\right)}𝐘_{\mathrm{𝟏𝟐𝟔}}^{\left(1\right)}`$ (4.37)
$`+(\stackrel{~}{\upsilon }_u^{\left(1\right)}3\stackrel{~}{\omega }_u^{\left(1\right)})𝐘_{\mathrm{𝟏𝟐𝟎}}^{\left(1\right)}+(\stackrel{~}{\upsilon }_u^{\left(2\right)}3\stackrel{~}{\omega }_u^{\left(2\right)})𝐘_{\mathrm{𝟏𝟐𝟎}}^{\left(2\right)}`$
$`𝐌_{\nu R}^{\left(\text{Maj}\right)}`$ $``$ $`M_I𝐘_{\mathrm{𝟏𝟐𝟔}}^{\left(1\right)}M_R𝐘_{\mathrm{𝟏𝟐𝟔}}^{\left(1\right)}`$ (4.38)
$`𝐌_{\nu L}^{\left(\text{Maj}\right)}`$ $``$ $`{\displaystyle \frac{\omega _u^{\left(1\right)2}}{M_I}}𝐘_{\mathrm{𝟏𝟐𝟔}}^{\left(1\right)}M_L𝐘_{\mathrm{𝟏𝟐𝟔}}^{\left(1\right)};M_L\omega _u^{\left(1\right)}`$ (4.39)
Für die einzelnen nichtverschwindenden Matrixelemente gilt unter Verwendung von (4.33) dann:
$`(𝐌_d)_{12}`$ $`=`$ $`x_1\upsilon _d^{\left(1\right)}+y_1\omega _d^{\left(1\right)}+z_1(\stackrel{~}{\upsilon }_d^{\left(1\right)}+\stackrel{~}{\omega }_d^{\left(1\right)})`$ (4.40)
$`(𝐌_d)_{21}`$ $`=`$ $`x_1\upsilon _d^{\left(1\right)}+y_1\omega _d^{\left(1\right)}z_1(\stackrel{~}{\upsilon }_d^{\left(1\right)}+\stackrel{~}{\omega }_d^{\left(1\right)})`$ (4.41)
$`(𝐌_d)_{23}`$ $`=`$ $`x_2\upsilon _d^{\left(2\right)}+z_2(\stackrel{~}{\upsilon }_d^{\left(2\right)}+\stackrel{~}{\omega }_d^{\left(2\right)})`$ (4.42)
$`(𝐌_d)_{32}`$ $`=`$ $`x_2\upsilon _d^{\left(2\right)}z_2(\stackrel{~}{\upsilon }_d^{\left(2\right)}+\stackrel{~}{\omega }_d^{\left(2\right)})`$ (4.43)
$`(𝐌_d)_{33}`$ $`=`$ $`\stackrel{~}{x}_1\upsilon _d^{\left(1\right)}+\stackrel{~}{y}_1\omega _d^{\left(1\right)}`$ (4.44)
$`=`$ $`\left({\displaystyle \frac{\stackrel{~}{x}_1}{x_1}}\right)x_1\upsilon _d^{\left(1\right)}+\left({\displaystyle \frac{\stackrel{~}{y}_1}{y_1}}\right)y_1\omega _d^{\left(1\right)}`$
$`(𝐌_e)_{12}`$ $`=`$ $`x_1\upsilon _d^{\left(1\right)}3y_1\omega _d^{\left(1\right)}+z_1(\stackrel{~}{\upsilon }_d^{\left(1\right)}3\stackrel{~}{\omega }_d^{\left(1\right)})`$ (4.45)
$`=`$ $`(𝐌_d)_{12}4y_1\omega _d^{\left(1\right)}4z_1\stackrel{~}{\omega }_d^{\left(1\right)}`$
$`(𝐌_e)_{21}`$ $`=`$ $`x_1\upsilon _d^{\left(1\right)}3y_1\omega _d^{\left(1\right)}z_1(\stackrel{~}{\upsilon }_d^{\left(1\right)}3\stackrel{~}{\omega }_d^{\left(1\right)})`$ (4.46)
$`=`$ $`(𝐌_d)_{21}4y_1\omega _d^{\left(1\right)}+4z_1\stackrel{~}{\omega }_d^{\left(1\right)}`$
$`(𝐌_e)_{23}`$ $`=`$ $`x_2\upsilon _d^{\left(2\right)}+z_2(\stackrel{~}{\upsilon }_d^{\left(2\right)}3\stackrel{~}{\omega }_d^{\left(2\right)})`$ (4.47)
$`=`$ $`(𝐌_d)_{23}4z_2\stackrel{~}{\omega }_d^{\left(2\right)}`$
$`(𝐌_e)_{32}`$ $`=`$ $`x_2\upsilon _d^{\left(2\right)}z_2(\stackrel{~}{\upsilon }_d^{\left(2\right)}3\stackrel{~}{\omega }_d^{\left(2\right)})`$ (4.48)
$`=`$ $`(𝐌_d)_{32}+4z_2\stackrel{~}{\omega }_d^{\left(2\right)}`$
$`(𝐌_e)_{33}`$ $`=`$ $`\stackrel{~}{x}_1\upsilon _d^{\left(1\right)}3\stackrel{~}{y}_1\omega _d^{\left(1\right)}`$ (4.49)
$`=`$ $`(𝐌_d)_{33}4\left({\displaystyle \frac{\stackrel{~}{y}_1}{y_1}}\right)y_1\omega _d^{\left(1\right)}`$
$`(𝐌_u)_{12}`$ $`=`$ $`x_1\upsilon _u^{\left(1\right)}+y_1\omega _u^{\left(1\right)}+z_1(\stackrel{~}{\upsilon }_u^{\left(1\right)}+\stackrel{~}{\omega }_u^{\left(1\right)})`$ (4.50)
$`(𝐌_u)_{21}`$ $`=`$ $`x_1\upsilon _u^{\left(1\right)}+y_1\omega _u^{\left(1\right)}z_1(\stackrel{~}{\upsilon }_u^{\left(1\right)}+\stackrel{~}{\omega }_u^{\left(1\right)})`$ (4.51)
$`(𝐌_u)_{23}`$ $`=`$ $`x_2\upsilon _u^{\left(2\right)}+z_2(\stackrel{~}{\upsilon }_u^{\left(2\right)}+\stackrel{~}{\omega }_u^{\left(2\right)})`$ (4.52)
$`(𝐌_u)_{32}`$ $`=`$ $`x_2\upsilon _u^{\left(2\right)}z_2(\stackrel{~}{\upsilon }_u^{\left(2\right)}+\stackrel{~}{\omega }_u^{\left(2\right)})`$ (4.53)
$`(𝐌_u)_{33}`$ $`=`$ $`\stackrel{~}{x}_1\upsilon _u^{\left(1\right)}+\stackrel{~}{y}_1\omega _u^{\left(1\right)}`$ (4.54)
$`=`$ $`\left({\displaystyle \frac{\stackrel{~}{x}_1}{x_1}}\right)x_1\upsilon _u^{\left(1\right)}+\left({\displaystyle \frac{\stackrel{~}{y}_1}{y_1}}\right)y_1\omega _u^{\left(1\right)}`$
$`(𝐌_\nu ^{\left(\text{Dir}\right)})_{12}`$ $`=`$ $`x_1\upsilon _u^{\left(1\right)}3y_1\omega _u^{\left(1\right)}+z_1(\stackrel{~}{\upsilon }_u^{\left(1\right)}3\stackrel{~}{\omega }_u^{\left(1\right)})`$ (4.55)
$`=`$ $`(𝐌_u)_{12}4y_1\omega _u^{\left(1\right)}4z_1\stackrel{~}{\omega }_u^{\left(1\right)}`$
$`(𝐌_\nu ^{\left(\text{Dir}\right)})_{21}`$ $`=`$ $`x_1\upsilon _u^{\left(1\right)}3y_1\omega _u^{\left(1\right)}z_1(\stackrel{~}{\upsilon }_u^{\left(1\right)}3\stackrel{~}{\omega }_u^{\left(1\right)})`$ (4.56)
$`=`$ $`(𝐌_u)_{21}4y_1\omega _u^{\left(1\right)}+4z_1\stackrel{~}{\omega }_u^{\left(1\right)}`$
$`(𝐌_\nu ^{\left(\text{Dir}\right)})_{23}`$ $`=`$ $`x_2\upsilon _u^{\left(2\right)}+z_2(\stackrel{~}{\upsilon }_u^{\left(2\right)}3\stackrel{~}{\omega }_u^{\left(2\right)})`$ (4.57)
$`=`$ $`(𝐌_u)_{23}4z_2\stackrel{~}{\omega }_u^{\left(2\right)}`$
$`(𝐌_\nu ^{\left(\text{Dir}\right)})_{32}`$ $`=`$ $`x_2\upsilon _u^{\left(2\right)}z_2(\stackrel{~}{\upsilon }_u^{\left(2\right)}3\stackrel{~}{\omega }_u^{\left(2\right)})`$ (4.58)
$`=`$ $`(𝐌_u)_{32}+4z_2\stackrel{~}{\omega }_u^{\left(2\right)}`$
$`(𝐌_\nu ^{\left(\text{Dir}\right)})_{33}`$ $`=`$ $`\stackrel{~}{x}_1\upsilon _u^{\left(1\right)}3\stackrel{~}{y}_1\omega _u^{\left(1\right)}`$ (4.59)
$`=`$ $`(𝐌_u)_{33}4\left({\displaystyle \frac{\stackrel{~}{y}_1}{y_1}}\right)y_1\omega _u^{\left(1\right)}`$
Eine der $`(\mathbf{15,2,2})_{120/126}`$ nicht an der Massenerzeugung teilnehmen zu lassen ist gleichbedeutend mit der Tatsache, daß entweder $`\omega _{u,d}^{\left(1\right)}`$, $`\stackrel{~}{\omega }_{u,d}^{\left(1\right)}`$ oder $`\stackrel{~}{\omega }_{u,d}^{\left(2\right)}`$ Null sein müssen. Wenn man bei fester Parameterzahl zunächst die größtmögliche Freiheit bezüglich der Gestalt der einzelnen Matrizen haben möchte, nämlich für alle $`j,k`$ $`(𝐌_d)_{jk}(𝐌_e)_{jk}`$ und $`(𝐌_u)_{jk}(𝐌_\nu ^{\left(\text{Dir}\right)})_{jk}`$ wählen zu können, führt das wegen der Gleichungen (4.45-4.49) und (4.55-4.59) zwangsläufig zu der Wahl $`\stackrel{~}{\omega }_{u,d}^{\left(1\right)}0`$.
Damit ist der Ansatz für die Massenmatrizen der Fermionen vollständig festgelegt. Alle Matrixeinträge setzen sich aus 14 voneinander unabhängigen freien $`SO(10)`$-Higgs-Parametern in Form von Produkten aus Vakuumerwartungswerten und Yukawa-Kopplungen zusammen. Dazu kommt noch die Massenskala $`M_R`$ der schweren Neutrinos, welche über den See-Saw-Mechanismus (2.51) auch als Vorfaktor $`(1/M_R)`$ in der See-Saw-Matrix erscheint und somit die Beträge der leichten Neutrinomassen festlegt.
Dem stehen 18 zu reproduzierende observable Größen gegenüber, nämlich die 12 Quark- und Leptonmassen sowie jeweils drei Mischungswinkel in V und U. Meßbar sind, zumindest bisher, nicht die eigentlichen Massen der leichten Neutrinos, sondern über die Oszillationsexperimente nur die beiden Parameter $`\mathrm{\Delta }m_{32}^2`$ und $`\mathrm{\Delta }m_{21}^2`$ (siehe Kapitel 3). Dem entspricht in gewisser Weise die Tatsache, daß der Wert von $`M_R`$ im Rahmen von GUT-Massenmodellen mit See-Saw-Mechanismus aufgrund des zu komplizierten vollständigen Higgs-Potentials nicht exakt berechenbar ist. Wie schon in Abschnitt 2.3.3 erläutert wurde, läßt sich lediglich die qualitative Relation $`M_RM_I`$ angeben, der genaue Wert von $`M_R`$ muß an die experimentellen Resultate angepaßt werden.
Es können demnach drei Vorhersagen gemacht werden, die im vorliegenden Fall alle im Neutrinosektor liegen. Allein aufgrund der Parameterzahl ist eine mögliche Erklärung der Anomalien von Sonnen- und atmosphärischen Neutrinos gemäß Abschnitt 3.2.4 als absolut nichttrivial zu betrachten. Die Zahl der freien Higgs-Parameter mag im Vergleich zu älteren Massenmodellen groß erscheinen; dafür werden in dem hier vorgeschlagenen Modell aber auch alle Massen der geladenen Fermionen sowie die CKM-Mischungswinkel bis auf Rundungsfehler korrekt wiedergegeben, während bisher Resultate $`\pm 10\%`$ vom tatsächlichen Wert schon als erfolgreich galten. Ferner wird zur Massenerzeugung auf keine Mechanismen zurückgegriffen, die, wie zum Beispiel nichtrenormierbare Operatoren, außerhalb der zugrundegelegten $`SO(10)`$-Theorie liegen.
### 4.4 Numerische Lösung des Massenmodells
Ausgangspunkt für den ersten Schritt in der numerischen Bestimmung der Massenmatrizen und Mischungen ist das Gleichungssystem
$`𝐔_L^{}𝐌_u𝐔_R`$ $`=`$ $`𝐌_u^{\left(D\right)},𝐃_L^{}𝐌_d𝐃_R=𝐌_d^{\left(D\right)},`$
$`𝐄_L^{}𝐌_e𝐄_R`$ $`=`$ $`𝐌_e^{\left(D\right)},𝐔_L^{}𝐃_L=𝐕,`$ (4.60)
welches, um die $`SO(10)`$-Beziehungen (4.34-4.39) ausnutzen zu können, bei $`\mu =M_I`$ gelöst werden muß. Auf den rechten Seiten der Gleichungen stehen dann die physikalischen Fermionmassen und die CKM-Matrix bei $`M_I`$. Bekannt sind deren Werte zunächst nur bei $`M_Z`$ (siehe (1.33) und Tabelle 1.6). Um mit Hilfe der Renormierungsgruppengleichungen des SM aus Anhang B.1 trotz Unkenntnis der Massenmatrizen bei $`M_Z`$ die entsprechenden Werte bei $`M_I`$ zu erhalten, werden zwei vereinfachende Annahmen gemacht:
* Die Skalenabhängigkeit der CKM-Matrix V ist sehr klein. Das ist in bestätigt worden.
* Die Teilchenmassen bei $`M_I`$ sind in guter Näherung unabhängig von der expliziten Form der Massenmatrizen bei $`M_Z`$. Das ist zumindest dann korrekt, wenn $`|(𝐌_u)_{33}|`$ sehr viel größer als alle anderen Einträge in $`𝐌_{u,d,e}`$ ist.
Inwieweit obige Annahmen im vorliegenden Fall berechtigt sind, wird sich später zeigen, wenn die Renormierungsgruppengleichungen der dann bekannten Yukawa-Matrizen von $`M_I`$ nach $`M_Z`$ integriert und diese bei $`M_Z`$ diagonalisiert werden.
Nun kann man die Yukawa-Matrizen bei $`M_Z`$ diagonal ansetzen und die Gleichungen (B.7-B.17) nach $`M_I=6.1410^{10}`$ GeV integrieren; die Resultate sind in Tabelle 4.1 zusammengefaßt.
In die Berechnungen geht auch die experimentell nicht bekannte Higgs-Selbstkopplung $`\lambda `$ ein, welche in niedrigster Ordnung über
$$m_H^2=\lambda \upsilon ^2;\upsilon =\mathrm{\hspace{0.33em}246.2}\text{GeV}$$
(4.61)
mit der Masse $`m_H`$ des SM-Higgs-Bosons zusammenhängt. Es muß also eine Annahme für $`\lambda (M_Z)`$ gemacht werden, wobei sich zeigt, daß $`\lambda (M_Z)=0.5`$ ein sinnvoller Wert zu sein scheint. Für $`\lambda (M_Z)0.4`$ führt die Integration von (B.7-B.17) nämlich auf $`\lambda (M_I)<0`$, während für $`\lambda (M_Z)0.6`$ der Wert von $`\lambda (M_I)`$ größer als 1 wird, was die Anwendbarkeit der Störungsrechnung gefährdet; dieses Phänomen ist in erwähnt worden. $`\lambda (M_Z)=0.5`$ würde in führender Ordnung einer Higgs-Masse von $`m_H174`$ GeV entsprechen, was mit den experimentellen Grenzen in vereinbar ist.
Damit sind die Matrizen auf den rechten Seiten von (4.60) bestimmt. Auf den linken Seiten stehen 31 unbekannte Größen:
* In den drei Massenmatrizen $`𝐌_f`$ ($`f=u,d,e`$) die 14 $`SO(10)`$-Higgs-Parameter aus (4.40-4.54). Es sind jedoch nur 13 voneinander unabhängig, da zwei lediglich in der Kombination $`z_2(\stackrel{~}{\upsilon }_u^{\left(2\right)}+\stackrel{~}{\omega }_u^{\left(2\right)})`$ vorkommen; um $`z_2\stackrel{~}{\upsilon }_u^{\left(2\right)}`$ und $`z_2\stackrel{~}{\omega }_u^{\left(2\right)}`$ separat bestimmen zu können, muß die Dirac-Massenmatrix der Neutrinos bekannt sein. Dies wird aber erst nach dem nächsten Schritt der Fall sein.
* Gemäß der Parametrisierung (1.34) jeweils drei Winkel $`\theta _{L,Rij}^{\left(f\right)}`$ ($`ij=\mathrm{12,23,31}`$) in den sechs Mischungsmatrizen, also insgesamt 18 Mischungswinkel.
Dem stehen 30 nichtlineare Bestimmungsgleichungen gegenüber, jeweils neun aus den ersten drei Beziehungen von (4.60) und drei unabhängige Gleichungen aus der letzten. Um das Problem numerisch behandeln zu können, muß also einer der freien Parameter vorgegeben werden. Hierfür wird der Quotient $`\stackrel{~}{y}_1/y_1`$ ausgewählt, welcher die Form der Majorana-Matrix (4.38) bis auf einen Vorfaktor eindeutig festlegt:
$$𝐌_{\nu R}^{\left(\text{Maj}\right)}=M_R\left(\begin{array}{ccc}0& y_1& 0\\ y_1& 0& 0\\ 0& 0& \stackrel{~}{y}_1\end{array}\right)=y_1M_R\left(\begin{array}{ccc}0& 1& 0\\ 1& 0& 0\\ 0& 0& \stackrel{~}{y}_1/y_1\end{array}\right)$$
(4.62)
Als möglicher Wertebereich für $`\stackrel{~}{y}_1/y_1`$ erscheint aufgrund der Massenhierarchie der Fermionen $`1|\stackrel{~}{y}_1/y_1|1000`$ plausibel. Bevor jedoch die Lösungen und ihre Eigenschaften behandelt werden, sei an dieser Stelle auf die Bedeutung der $`SO(10)`$-spezifischen Beziehungen (4.34-4.37) hingewiesen.
Wie man an (4.40-4.49) erkennt, sind bei bekanntem $`𝐌_d`$ in $`𝐌_e`$ noch drei Freiheitsgrade enthalten, nämlich $`y_1\omega _d^{\left(1\right)}`$, $`\stackrel{~}{y}_1/y_1`$ und $`z_2\stackrel{~}{\omega }_d^{\left(2\right)}`$. Da die Leptonmassen sich auch bei $`M_I`$ deutlich von den Massen der $`d`$-Quarks unterscheiden, müssen diese Parameter derart festgelegt werden, daß $`𝐌_e`$ die korrekten Teilchenmassen liefert. Dann sind die $`SO(10)`$-Higgs-Parameter im $`d`$-$`e`$-Sektor eindeutig bestimmt:
$`x_1\upsilon _d^{\left(1\right)}`$ $`=`$ $`{\displaystyle \frac{1}{8}}\left(\mathrm{\hspace{0.17em}3}(𝐌_d)_{12}+3(𝐌_d)_{21}+(𝐌_e)_{12}+(𝐌_e)_{21}\right)`$ (4.63)
$`y_1\omega _d^{\left(1\right)}`$ $`=`$ $`{\displaystyle \frac{1}{8}}\left((𝐌_d)_{12}+(𝐌_d)_{21}(𝐌_e)_{12}(𝐌_e)_{21}\right)`$ (4.64)
$`z_1\stackrel{~}{\upsilon }_d^{\left(1\right)}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left((𝐌_d)_{12}(𝐌_d)_{21}\right)={\displaystyle \frac{1}{2}}\left((𝐌_e)_{12}(𝐌_e)_{21}\right)`$ (4.65)
$`\stackrel{~}{x}_1\upsilon _d^{\left(1\right)}`$ $`=`$ $`{\displaystyle \frac{1}{4}}\left(\mathrm{\hspace{0.17em}3}(𝐌_d)_{33}+(𝐌_e)_{33}\right)`$ (4.66)
$`\stackrel{~}{y}_1\omega _d^{\left(1\right)}`$ $`=`$ $`{\displaystyle \frac{1}{4}}\left((𝐌_d)_{33}(𝐌_e)_{33}\right)`$ (4.67)
$`x_2\upsilon _d^{\left(2\right)}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left((𝐌_d)_{23}+(𝐌_d)_{32}\right)={\displaystyle \frac{1}{2}}\left((𝐌_e)_{23}+(𝐌_e)_{32}\right)`$ (4.68)
$`z_2\stackrel{~}{\upsilon }_d^{\left(2\right)}`$ $`=`$ $`{\displaystyle \frac{1}{8}}\left(\mathrm{\hspace{0.17em}3}(𝐌_d)_{23}3(𝐌_d)_{32}+(𝐌_e)_{23}(𝐌_e)_{32}\right)`$ (4.69)
$`z_2\stackrel{~}{\omega }_d^{\left(2\right)}`$ $`=`$ $`{\displaystyle \frac{1}{8}}\left((𝐌_d)_{23}(𝐌_d)_{32}(𝐌_e)_{23}+(𝐌_e)_{32}\right)`$ (4.70)
$`{\displaystyle \frac{\stackrel{~}{x}_1}{x_1}}`$ $`=`$ $`2{\displaystyle \frac{3(𝐌_d)_{33}+(𝐌_e)_{33}}{3(𝐌_d)_{12}+3(𝐌_d)_{21}+(𝐌_e)_{12}+(𝐌_e)_{21}}}`$ (4.71)
$`{\displaystyle \frac{\stackrel{~}{y}_1}{y_1}}`$ $`=`$ $`2{\displaystyle \frac{(𝐌_d)_{33}(𝐌_e)_{33}}{(𝐌_d)_{12}+(𝐌_d)_{21}(𝐌_e)_{12}(𝐌_e)_{21}}}`$ (4.72)
$`{\displaystyle \frac{\stackrel{~}{\upsilon }_d^{\left(1\right)}}{\stackrel{~}{\omega }_d^{\left(1\right)}}}`$ $`=`$ $`{\displaystyle \frac{3(𝐌_d)_{12}3(𝐌_d)_{21}+(𝐌_e)_{12}(𝐌_e)_{21}}{(𝐌_d)_{12}(𝐌_d)_{21}(𝐌_e)_{12}+(𝐌_e)_{21}}}`$ (4.73)
$`{\displaystyle \frac{\stackrel{~}{\upsilon }_d^{\left(2\right)}}{\stackrel{~}{\omega }_d^{\left(2\right)}}}`$ $`=`$ $`{\displaystyle \frac{3(𝐌_d)_{23}3(𝐌_d)_{32}+(𝐌_e)_{23}(𝐌_e)_{32}}{(𝐌_d)_{23}(𝐌_d)_{32}(𝐌_e)_{23}+(𝐌_e)_{32}}}`$ (4.74)
Ist nun auch $`𝐌_u`$ bekannt, enthält die Dirac-Massenmatrix der Neutrinos $`𝐌_\nu ^{\left(\text{Dir}\right)}`$ gemäß (4.50-4.59) relativ zu $`𝐌_u`$ nur noch den Freiheitsgrad $`z_2\stackrel{~}{\omega }_u^{\left(2\right)}`$, da $`\stackrel{~}{y}_1/y_1`$ aus dem $`d`$-$`e`$-Sektor bestimmt ist. Auch $`y_1\omega _u^{\left(1\right)}`$ ist durch die Kenntnis von $`𝐌_u`$ bereits eindeutig festgelegt, wie die zweite der folgenden Gleichungen zeigt:
$`x_1\upsilon _u^{\left(1\right)}`$ $`=`$ $`{\displaystyle \frac{{\displaystyle \frac{\stackrel{~}{y}_1}{2y_1}}\left((𝐌_u)_{12}+(𝐌_u)_{21}\right)+(𝐌_u)_{33}}{{\displaystyle \frac{\stackrel{~}{x}_1}{x_1}}{\displaystyle \frac{\stackrel{~}{y}_1}{y_1}}}}`$ (4.75)
$`y_1\omega _u^{\left(1\right)}`$ $`=`$ $`{\displaystyle \frac{{\displaystyle \frac{\stackrel{~}{x}_1}{2x_1}}\left((𝐌_u)_{12}+(𝐌_u)_{21}\right)(𝐌_u)_{33}}{{\displaystyle \frac{\stackrel{~}{x}_1}{x_1}}{\displaystyle \frac{\stackrel{~}{y}_1}{y_1}}}}`$ (4.76)
$`z_1\stackrel{~}{\upsilon }_u^{\left(1\right)}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left((𝐌_u)_{12}(𝐌_u)_{21}\right)={\displaystyle \frac{1}{2}}\left((𝐌_\nu ^{\left(\text{Dir}\right)})_{12}(𝐌_\nu ^{\left(\text{Dir}\right)})_{21}\right)`$ (4.77)
$`\stackrel{~}{x}_1\upsilon _u^{\left(1\right)}`$ $`=`$ $`\left({\displaystyle \frac{\stackrel{~}{x}_1}{x_1}}\right)x_1\upsilon _u^{\left(1\right)}`$ (4.78)
$`\stackrel{~}{y}_1\omega _u^{\left(1\right)}`$ $`=`$ $`\left({\displaystyle \frac{\stackrel{~}{y}_1}{y_1}}\right)y_1\omega _u^{\left(1\right)}`$ (4.79)
$`x_2\upsilon _u^{\left(2\right)}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left((𝐌_u)_{23}+(𝐌_u)_{32}\right)={\displaystyle \frac{1}{2}}\left((𝐌_\nu ^{\left(\text{Dir}\right)})_{23}+(𝐌_\nu ^{\left(\text{Dir}\right)})_{32}\right)`$ (4.80)
$`z_2\stackrel{~}{\upsilon }_u^{\left(2\right)}`$ $`=`$ $`{\displaystyle \frac{1}{8}}\left(\mathrm{\hspace{0.17em}3}(𝐌_u)_{23}3(𝐌_u)_{32}+(𝐌_\nu ^{\left(\text{Dir}\right)})_{23}(𝐌_\nu ^{\left(\text{Dir}\right)})_{32}\right)`$ (4.81)
$`z_2\stackrel{~}{\omega }_u^{\left(2\right)}`$ $`=`$ $`{\displaystyle \frac{1}{8}}\left((𝐌_u)_{23}(𝐌_u)_{32}(𝐌_\nu ^{\left(\text{Dir}\right)})_{23}+(𝐌_\nu ^{\left(\text{Dir}\right)})_{32}\right)`$ (4.82)
Anders ausgedrückt, wenn man einen Wert für $`\stackrel{~}{y}_1/y_1`$ vorgibt und Lösungen für die Massenmatrizen der geladenen Fermionen findet, welche (4.60) erfüllen, so ist der Neutrinosektor der Theorie bis auf die beiden Parameter $`y_1M_R`$ und $`z_2\stackrel{~}{\omega }_u^{\left(2\right)}`$ eindeutig festgelegt. Das wiederum deutet schon den zweiten Schritt der Analyse an, nämlich die Suche nach Werten für $`\stackrel{~}{y}_1/y_1`$ und $`z_2\stackrel{~}{\omega }_u^{\left(2\right)}`$, welche eine Oszillationslösung für die Anomalien der Sonnen- und atmosphärischen Neutrinos gemäß 3.2.4 ermöglichen.
Die Suche nach Lösungen von (4.60) für vorgegebene Werte von $`\stackrel{~}{y}_1/y_1`$ mit $`1|\stackrel{~}{y}_1/y_1|1000`$ liefert folgende Resultate:
* Für $`|\stackrel{~}{y}_1/y_1|5`$ und $`|\stackrel{~}{y}_1/y_1|500`$ können keine Lösungen gefunden werden.
* Für $`5|\stackrel{~}{y}_1/y_1|500`$ existieren Lösungen. Es gibt bei gegebenem $`\stackrel{~}{y}_1/y_1`$ bis zu sechs (bei $`20\stackrel{~}{y}_1/y_125`$) verschiedene Sätze von Quark-Massenmatrizen und zu jedem dieser Paare jeweils zwei Möglichkeiten für die Massenmatrix der geladenen Leptonen, welche (4.60) erfüllen. Letzteres resultiert aus der Tatsache, daß es zu jeder Matrix $`𝐌_d`$ zwei Werte des Parameters $`z_2\stackrel{~}{\omega }_d^{\left(2\right)}`$ gibt, welche die korrekten Leptonmassen liefern; der Parameter $`y_1\omega _d^{\left(1\right)}`$ dagegen ist in beiden Fällen gleich groß.
Es werden nur ganzzahlige Werte für $`\stackrel{~}{y}_1/y_1`$ untersucht, da die Lösungen, sofern sie existieren, stetig von $`\stackrel{~}{y}_1/y_1`$ abhängen. Man erhält mit dem Ansatz (4.34-4.39) für die Massenmatrizen demnach eine Vielzahl von Lösungen für die Matrizen der geladenen Fermionen, welche die beobachteten Massen und CKM-Mischungen liefern. Nun muß der Neutrinosektor in die Analyse mit einbezogen werden.
Es verbleiben zwei Parameter, welche aus den nun bekannten Massenmatrizen der geladenen Fermionen gemäß den Beziehungen (4.63-4.82) nicht berechnet werden können, nämlich $`y_1M_R`$ und $`z_2\stackrel{~}{\omega }_u^{\left(2\right)}`$. Ersterer kann im Rahmen des Massenmodells nicht exakt vorhergesagt werden, sondern muß am Ende an die Resultate der Oszillationsexperimente angepaßt werden, wobei $`M_RM_I`$ gelten sollte. Bei vorgegebenem $`\stackrel{~}{y}_1/y_1`$ ist demnach $`z_2\stackrel{~}{\omega }_u^{\left(2\right)}`$ der einzige noch freie Parameter in der Dirac-Massenmatrix der Neutrinos; die Majorana-Matrix ist bis auf den globalen Vorfaktor $`y_1M_R`$ festgelegt. Die Aufgabe besteht also darin, zu untersuchen, ob für diejenigen Werte von $`\stackrel{~}{y}_1/y_1`$, welche Lösungen von (4.60) liefern, der Parameter $`z_2\stackrel{~}{\omega }_u^{\left(2\right)}`$ so gewählt werden kann, daß die Massen und Mischungen der leichten Neutrinos phänomenologisch sinnvoll sind, das heißt den experimentellen Grenzen in (3.20-3.33) genügen. Da das Fermionspektrum des Modells drei Neutrino-Arten enthält, können durch Oszillationslösungen zwei der drei in Kapitel 3 diskutierten Neutrinoprobleme erklärt werden; das sollen aufgrund der im Vergleich zu den LSND-Resultaten deutlich stärkeren Evidenz die Anomalien der Sonnen- und atmosphärischen Neutrinos sein. Wegen der Argumentation in Abschnitt 3.2.1 wird für das Sonnen-Neutrinodefizit nur der MSW-Effekt berücksichtigt.
Konkret sieht die Vorgehensweise wie folgt aus:
* Für festes $`\stackrel{~}{y}_1/y_1`$ wird der gesamte plausibel erscheinende Wertebereich für $`z_2\stackrel{~}{\omega }_u^{\left(2\right)}`$ untersucht. Im folgenden wird $`|z_2\stackrel{~}{\omega }_u^{\left(2\right)}|m_t(M_I)`$ gewählt.
* Bei gegebenem $`z_2\stackrel{~}{\omega }_u^{\left(2\right)}`$-Wert werden die Dirac-, Majorana- und See-Saw-Massenmatrizen der Neutrinos bestimmt; letztere wird gemäß (2.53) diagonalisiert.
* Auf diese Weise erhält man die Massen $`m_{\nu _i}`$ der leichten Neutrinos bis auf den Vorfaktor $`1/(y_1M_R)`$, es wird $`|m_{\nu _3}||m_{\nu _2}||m_{\nu _1}|`$ sowie $`\mathrm{\Delta }m_{\text{atm}}^2m_{\nu _3}^2m_{\nu _2}^2`$ und $`\mathrm{\Delta }m_{\text{sun}}^2m_{\nu _2}^2m_{\nu _1}^2`$ angenommen. Die leptonische Mischungsmatrix $`𝐔𝐄_L^{}𝐍_L`$ ist damit ebenfalls bekannt.
* Nun wird überprüft, ob die so erhaltenen Neutrinoeigenschaften den Einschränkungen (3.20-3.33) genügen. In Unkenntnis der Massenskala $`M_R`$ werden folgende Forderungen an die Lösungen gestellt:
$`|(𝐔)_{13}|0.05`$ (4.83)
$`0.49|(𝐔)_{23}|0.71`$ (4.84)
$`0.03|(𝐔)_{12}|0.05\text{(MSW mit kleiner Mischung)}`$
$`\text{oder}0.35|(𝐔)_{12}|0.49\text{(MSW mit großer Mischung)}`$ (4.85)
$`50\mathrm{\Delta }m_{\text{atm}}^2/\mathrm{\Delta }m_{\text{sun}}^2(m_{\nu _3}^2m_{\nu _2}^2)/(m_{\nu _2}^2m_{\nu _1}^2)1000`$ (4.86)
Es zeigt sich, daß drei Bereiche im $`\stackrel{~}{y}_1/y_1`$-$`z_2\stackrel{~}{\omega }_u^{\left(2\right)}`$-Parameterraum existieren, in denen Lösungen von (4.60) liegen, welche zudem im Neutrinosektor die Forderungen (4.83-4.86) erfüllen. Die Abbildungen 4.1, 4.2 und 4.3 zeigen die Lage der entsprechenden Parameterbereiche, wobei zu jedem untersuchten ganzzahligen $`\stackrel{~}{y}_1/y_1`$-Wert die Minimal- und Maximalwerte von $`z_2\stackrel{~}{\omega }_u^{\left(2\right)}`$ angegeben sind, die zu Neutrinoeigenschaften führen, welche den Einschränkungen (4.83-4.86) genügen. Ein Bereich liefert MSW-Lösungen für das Sonnen-Neutrinodefizit mit großem Mischungswinkel, die anderen beiden Bereiche liefern MSW-Lösungen mit kleiner Mischung.
Aus jedem der drei Bereiche wird nun eine repräsentative Lösung ausgewählt, die genauer untersucht werden soll. Deren Werte für $`\stackrel{~}{y}_1/y_1`$ und $`z_2\stackrel{~}{\omega }_u^{\left(2\right)}`$ werden so gewählt, daß die zugehörigen Punkte ungefähr in der Mitte der erlaubten Parameterbereiche liegen, da sowohl die bestehenden als auch die neu hinzukommenden Neutrinoexperimente die Grenzen in (3.20-3.33) weiter einschränken werden. Tabelle 4.2 faßt die Parameterwerte für die drei Beispielmodelle zusammen.
Die numerischen Resultate für die untersuchten Beispielmodelle sind in Anhang D zusammengefaßt. Zunächst gibt Abschnitt D.1 Auskunft über die nichtverschwindenden Einträge aller Massenmatrizen sowie die Fermionmassen und -mischungen bei $`M_I`$. Daraus lassen sich mit Hilfe von (4.63-4.82) die Werte der $`SO(10)`$-Higgs-Parameter berechnen; die Resultate enthält Abschnitt D.2. Um die Werte dieser Größen bei $`M_Z`$ zu erhalten, müssen die Renormierungsgruppengleichungen (B.7-B.20) des SM einschließlich derjenigen für die See-Saw-Massenmatrix von $`M_I`$ nach $`M_Z`$ integriert werden. Dabei werden auch die Matrixeinträge ungleich Null, die bei $`M_I`$ wegen des Ansatzes (4.31) verschwinden. Die Elemente der Massenmatrizen der geladenen Fermionen und leichten Neutrinos sowie deren Massen und Mischungen bei $`M_Z`$ sind in Abschnitt D.3 aufgelistet. Insbesondere die Mischungswinkel werden bei der Berechnung der Nukleonzerfallsraten im nächsten Kapitel von großer Bedeutung sein. Aus ihnen lassen sich auch die CKM-Matrix V und die leptonische Mischungsmatrix U bestimmen; die Ergebnisse für die drei betrachteten Modelle enthält Abschnitt D.4.
Obwohl die Massenmatrizen in den einzelnen Fällen durchaus verschieden sind, lassen sich folgende gemeinsame Eigenschaften der Lösungen feststellen:
* Die Winkel der CKM-Matrix in der hier für alle Mischungsmatrizen verwendeten Parametrisierung (1.34) sind mit $`\theta _{12}=0.223`$, $`\theta _{23}=0.039`$ und $`\theta _{31}=0.003`$ angesetzt worden, was innerhalb der experimentellen Grenzen liegt. Es zeigt sich, daß die linkshändigen Mischungswinkel der Quarks $`\theta _{Ljk}^{\left(u\right)}`$ und $`\theta _{Ljk}^{\left(d\right)}`$ betragsmäßig zum Teil deutlich größer als die CKM-Winkel sind. Besonders auffällig ist das im Fall von $`\theta _{12}\theta _{L12}^{\left(d\right)}\theta _{L12}^{\left(u\right)}`$. Demnach kann man von der Größe der observablen CKM-Winkel nicht zwangsläufig auf die der linkshändigen Quarkmischungen schließen, wie es häufig getan wird.
* Alle Modelle enthalten betragsmäßig große Mischungswinkel mit Beträgen zwischen $`0.5`$ und $`1.0`$; so ist zum Beispiel $`\theta _{R23}^{\left(d\right)}`$ in allen Fällen sehr groß. Große rechtshändige Mischungen sind insofern von besonderem Interesse, als sie im Rahmen des SM keine experimentell beobachtbaren Auswirkungen haben, in GUT-Modellen wie dem vorliegenden aber durchaus relevant sind. Unter anderem hängen die Verzweigungsraten der Nukleonenzerfälle stark von den Fermionmischungen ab, wie man im nächsten Kapitel sehen wird.
* In den Lösungen ist $`(𝐌_u)_{33}m_t`$ und somit betragsmäßig deutlich größer als alle übrigen Matrixeinträge, $`(𝐌_d)_{33}`$ und $`(𝐌_e)_{33}`$ dagegen sind in $`𝐌_d`$ beziehungsweise $`𝐌_e`$ nicht allein dominant. Die $`(𝐌_f)_{12}`$-$`(𝐌_f)_{21}`$\- und $`(𝐌_f)_{23}`$-$`(𝐌_f)_{32}`$-Asymmetrien in den Massenmatrizen sind in allen Fällen stark ausgeprägt.
* Die zu Beginn des Abschnitts gemachten Annahmen über die Skalenabhängigkeit der Fermionmassen und CKM-Mischungen sind gerechtfertigt. Alle Werte für die Massen bei $`M_Z`$, welche die drei Modelle durch Integration der Renormierungsgruppengleichungen liefern, ergeben sich abgesehen von Rundungsfehlern in sehr guter Übereinstimmung mit den ursprünglich angesetzten experimentellen Größen. Auch die aus den Lösungen numerisch bestimmten CKM-Matrizen genügen den experimentellen Grenzen. Lediglich die Skalenabhängigkeit von $`\theta _{L23}^{\left(d\right)}`$ und $`\theta _{L31}^{\left(d\right)}`$ ist nicht völlig zu vernachlässigen; es gilt $`\theta _{L23}^{\left(d\right)}(M_I)>\theta _{L23}^{\left(d\right)}(M_Z)`$ und $`\theta _{L31}^{\left(d\right)}(M_I)<\theta _{L31}^{\left(d\right)}(M_Z)`$. Das wiederum führt innerhalb der Grenzen zu vergleichsweise großen Werten für $`(𝐕)_{33}`$, während $`(𝐕)_{23}`$ und $`(𝐕)_{32}`$ betragsmäßig am unteren Ende des erlaubten Bereichs liegen. Man könnte diese Schwäche jedoch ohne weiteres beheben, indem man die nun bekannte Skalenabhängigkeit auf den rechten Seiten von (4.60) berücksichtigt. Da die Effekte aber sehr klein sind, wird hier darauf verzichtet.
Die drei untersuchten Beispiellösungen beschreiben den Sektor der geladenen Fermionen in sehr guter Übereinstimmung mit den experimentellen Resultaten für die Massen und Mischungen. Das ist jedoch lediglich ein notwendiges Kriterium dafür, daß der Ansatz (4.31) als sinnvoll zu betrachten ist. Auf die Neutrinoeigenschaften der einzelnen Lösungen, welche die eigentlichen Vorhersagen des Massenmodells darstellen, wird im nächsten Abschnitt eingegangen.
Die in diesem Kapitel verwendeten numerischen Algorithmen sind im wesentlichen entnommen. Zur Integration der Renormierungsgruppengleichungen wurde eine auf der Runge-Kutta-Methode vierter Ordnung basierende Routine mit adaptiver Schrittweitensteuerung benutzt. Dabei wird der durch die endliche Schrittweite erzeugte Fehler abgeschätzt und unterhalb einer vorzugebenden Schranke gehalten. Diese Schranke wurde so festgelegt, daß die Fehler aufgrund der Integration deutlich kleiner als die experimentellen Unsicherheiten sind. In den $`\beta `$-Funktionskoeffizienten zweiter Ordnung der Yukawa-Matrizen und der Higgs-Kopplung wurde die Näherung benutzt, nur die betragsmäßig größten Einträge der Yukawa-Matrizen zu verwenden und alle übrigen auf Null zu setzen.
Die Diagonalisierung der Massenmatrizen gemäß (1.29) entspricht der numerischen Lösung eines nichtlinearen Gleichungssystems und wurde mit Hilfe der multidimensionalen Newton-Raphson-Methode durchgeführt.
### 4.5 Eigenschaften der leichten Neutrinos
Die Massen und Mischungen der leichten Neutrinos bei $`M_Z`$ sind in Tabelle 4.3 zusammengefaßt, welche einen Ausschnitt aus Tabelle D.4 darstellt.
Da sich die Massen der geladenen Leptonen zwischen $`\mu =M_Z`$ und $`\mu /M_Z0`$ kaum ändern, kann man das in noch stärkerem Maße von den Neutrinomassen erwarten, da die Neutrinos nur schwach wechselwirken. Vergleicht man die Tabellen D.4 einerseits und D.1 sowie D.2 andererseits, so erkennt man, daß die Mischungswinkel aller Fermionen zwischen $`M_I`$ und $`M_Z`$ nur minimal skalenabhängig sind. Im folgenden wird deshalb von $`\theta _{L,Rjk}^{\left(f\right)}(M_Z)=\theta _{L,Rjk}^{\left(f\right)}(m_p)`$ ausgegangen.
In Tabelle 4.3 erkennt man zunächst die starke Ähnlichkeit der Eigenschaften von Modell 2a und 2b. Beide realisieren eine Oszillationslösung des Sonnen-Neutrinodefizits über den MSW-Effekt mit kleinen Mischungen, während Modell 1 auf den MSW-Effekt mit großer Mischung führt. Von den linkshändigen Mischungen der geladenen Leptonen ist lediglich $`\theta _{L23}^{\left(e\right)}`$ in Modell 1 betragsmäßig deutlich größer als Null. In den anderen beiden Fällen kommt die große (23)-Mischung in U, welche die Anomalie der atmosphärischen Neutrinos erklärt, allein durch die reine Neutrinomischung zustande.
Für $`(m_{\nu _3}^2m_{\nu _2}^2)/(m_{\nu _2}^2m_{\nu _1}^2)\mathrm{\Delta }m_{\text{atm}}^2/\mathrm{\Delta }m_{\text{sun}}^2`$ waren aufgrund der experimentellen Grenzen Werte zwischen 50 und 1000 zugelassen worden. Hier fällt auf, daß der Wert in Modell 1 nahe an der oberen Schranke ist, während die Resultate in den Modellen 2a,b deutlich im unteren Teil des erlaubten Bereichs liegen. Dem steht die durch Abbildung 3.1(a) begründete Erwartung gegenüber, bei festem $`\mathrm{\Delta }m_{\text{atm}}^2`$ für den MSW-Effekt mit großer Mischung den im Vergleich zum MSW-Effekt mit kleiner Mischung kleineren Wert für $`\mathrm{\Delta }m_{\text{atm}}^2/\mathrm{\Delta }m_{\text{sun}}^2`$ zu finden.
Tabelle 4.4 gibt die Werte von $`(y_1M_R/M_I)`$ und $`y_1M_R`$ an, welche benötigt werden, um die erhaltenen Vorhersagen für $`(m_{\nu _3}^2m_{\nu _2}^2)`$ und $`(m_{\nu _2}^2m_{\nu _1}^2)`$ an die experimentellen Grenzen für $`\mathrm{\Delta }m_{\text{atm}}^2`$ und $`\mathrm{\Delta }m_{\text{sun}}^2`$ aus (3.20-3.33) anzupassen.
Sie sind angesichts der erwarteten Relation $`M_RM_I`$ vergleichsweise groß, da $`y_1`$ als Yukawa-Kopplung $`1`$ sein sollte. Andererseits ist $`(y_1M_R/M_I)100`$ auch keineswegs ausgeschlossen, da der exakte Wert von $`M_R`$ von zahlreichen Parametern im nicht explizit bekannten Higgs-Potential abhängt. Diese sollten zwar von der Größenordnung 1 sein, aber es ist durchaus möglich, daß der Proportionalitätskoeffizient in $`M_RM_I`$ aus einer nichttrivialen Kombination der Parameter besteht und im Bereich $`1/100`$ bis $`100`$ liegt. Überdies ist in Abschnitt 4.2 eine Unsicherheit $`10^{\pm 2.6}400`$ in $`M_I`$ aufgrund von Schwelleneffekten berechnet worden. Die Vernachlässigung der massenabhängigen Schwellenkorrekturen bei der Bestimmung von $`M_I`$ in Abschnitt 4.1 könnte somit theoretisch zu einem um zwei Größenordnungen zu kleinen Wert geführt haben. Eine Kombination dieser beiden Möglichkeiten ist demnach definitiv in der Lage, $`(y_1M_R/M_I)100`$ zu realisieren.
Wie schon weiter oben erwähnt wurde, liegen im Modell 1 die Werte für $`m_{\nu _2}^2m_{\nu _1}^2`$ im unteren Teil des für $`\mathrm{\Delta }m_{\text{sun}}^2`$ zulässigen Bereichs, während jene für $`m_{\nu _3}^2m_{\nu _2}^2`$ im Rahmen der experimentellen Grenzen von $`\mathrm{\Delta }m_{\text{atm}}^2`$ relativ groß sind. In den Modellen 2a,b beobachtet man ein gegenteiliges Verhalten dieser Größen.
Vergleicht man die Modellvorhersagen für die leptonische Mischungsmatrix U (siehe Anhang D.4) mit den Grenzen aus (3.20-3.33), so erkennt man eine gute Übereinstimmung zwischen beiden; alle Einträge von U liegen innerhalb der erlaubten Bereiche. Auffällig ist bei allen drei Modellen die Tatsache, daß die (23)-Mischung in U vergleichsweise klein ist. Insbesondere in Modell 1 ist $`|(𝐔)_{23}|=0.504`$ nahe an der unteren Schranke 0.49, während die $`|(𝐔)_{23}|`$-Werte in den Modellen 2a,b ($`0.536`$ beziehungsweise $`0.552`$) etwas größer sind. Da die existierenden ebenso wie die neu hinzukommenden Neutrinoexperimente die Grenzen in (3.20-3.33) weiter einschränken werden, ist Modell 1 in diesem Sinne weniger überzeugend als die anderen beiden. Das hier vorgeschlagene $`SO(10)`$-Massenmodell deutet also auf eine MSW-Lösung mit kleiner Mischung hin. Sollte der MSW-Effekt mit kleiner Mischung durch Experimente eindeutig als Ursache für das Defizit der Sonnen-Neutrinos bestätigt werden, wird es schwierig sein, eines der Modelle 2a und 2b anhand der Eigenschaften im Neutrinosektor als in der Natur realisiert zu identifizieren. Dazu muß man auf die Vorhersagen bezüglich der Nukleonzerfallsraten zurückgreifen, welche im nächsten Kapitel berechnet werden.
Die Existenz von leichten Majorana-Neutrinos hat eine weitere überprüfbare Konsequenz, nämlich die sogenannten neutrinolosen doppelten $`\beta `$-Zerfälle. Sie kommen durch den Austausch eines Majorana-Neutrinos zwischen zwei zerfallenden Neutronen zustande, was zu Kernreaktionen der Form $`X_Z^AX_{Z+2}^A+2e^{}`$ führt. Die experimentelle Nichtbeobachtung solcher Prozesse liefert die Einschränkung :
$$m_\nu \underset{i}{}\left((𝐔)_{1i}^2m_{\nu _i}\right)\mathrm{\hspace{0.33em}0.2}\text{eV}$$
(4.87)
Alle hier untersuchten Lösungen erfüllen diese Bedingung.
### 4.6 Baryonasymmetrie durch Neutrinozerfälle
Obwohl nicht Gegenstand dieser Arbeit, soll hier kurz auf die Möglichkeit hingewiesen werden, die Baryonasymmetrie des Universums im Rahmen einer $`SO(10)`$-GUT mit See-Saw-Mechanismus durch den Zerfall der schweren Majorana-Neutrinos zu erklären. Einen aktuellen Überblick zu diesem Thema liefert .
Das beobachtete Universum scheint wesentlich mehr Materie als Antimaterie zu enthalten. Wenn $`n_b`$, $`n_{\overline{b}}`$ und $`n_\gamma `$ die mittleren Teilchenzahldichten der Baryonen, Antibaryonen und Photonen im Universum bezeichnen, ergibt sich für die Größe $`\eta (n_bn_{\overline{b}})/n_\gamma `$ aus kosmologischen Überlegungen der Wert $`410^{10}\eta 710^{10}`$ . Wenn man davon ausgeht, daß $`\eta `$ bei der Entstehung des Universums Null war, so kann die Kosmologie die heute beobachtete Baryonasymmetrie nicht erklären. Dazu muß auf die Physik der Elementarteilchen zurückgegriffen werden.
Wie in gezeigt wurde, sind für die Entstehung der Baryonasymmetrie in einem ursprünglich symmetrischen Universum drei notwendige Kriterien zu erfüllen, nämlich die Existenz baryonzahlverletzender Wechselwirkungen, die Verletzung von $`C`$ und $`CP`$ sowie das Vorhandensein eines thermischen Ungleichgewichts. Die Baryonzahlverletzung ist naheliegenderweise notwendig, um $`\eta 0`$ zu ermöglichen, und ist sowohl in $`SO(10)`$\- als auch $`G_{\text{PS}}`$-Modellen enthalten; lediglich die Kombination $`BL`$ ist Teil der Eichsymmetrie und somit eine Erhaltungsgröße. $`C`$\- und $`CP`$-Verletzung sind erforderlich, da sich zeigen läßt, daß andernfalls die Rate der Prozesse, welche Baryonen erzeugen, genauso groß wie die der antibaryonerzeugenden Prozesse ist. Während die Eichboson- und Fermionsektoren von $`SO(10)`$-GUTs $`C`$\- und $`CP`$-invariant sind, lassen sich diese Symmetrien durch Einführung komplexer Higgs-Kopplungen explizit brechen. Schließlich ist ein thermisches Ungleichgewicht erforderlich, da der Baryonzahloperator $`B`$ bei einer Temperatur $`T`$ im Gleichgewicht den Erwartungswert
$`B_T`$ $`=`$ $`\text{Tr}\left(e^{\beta H}B\right)=\text{Tr}\left([𝒞𝒫𝒯][𝒞𝒫𝒯]^1e^{\beta H}B\right)`$ (4.88)
$`=`$ $`\text{Tr}\left(e^{\beta H}[𝒞𝒫𝒯]^1B[𝒞𝒫𝒯]\right)=\text{Tr}\left(e^{\beta H}B\right)\stackrel{!}{=}\mathrm{\hspace{0.33em}0}`$
besitzt, vorausgesetzt, die Theorie ist $`CPT`$-invariant. Ob sich ein System im thermischen Gleichgewicht befindet, läßt sich nur durch Lösung seiner Boltzmann-Gleichungen feststellen. In erster Näherung kann man aber sagen, daß ein Ungleichgewicht vorliegt, wenn die Expansionsrate des Universums größer als die Raten der baryon- und leptonzahlverletzenden Wechselwirkungen ist. Das sollte bei Temperaturen $`TM_I`$ der Fall sein.
Somit ist im vorliegenden $`SO(10)`$-Modell die Entstehung einer $`B+L`$-Asymmetrie bei Energien $`\mu M_I`$ möglich. Allerdings existieren auch im SM baryon- und leptonzahlverletzende Prozesse, welche $`BL`$ erhalten und auf nichtperturbativen Effekten beruhen. Während nämlich die klassische Lagrangedichte des SM invariant unter globalen $`U(1)_{B,L}`$-Symmetrien ist, besitzen die zu diesen Symmetrien gehörenden Ströme in der quantisierten Theorie Anomalien . Das führt zu $`B+L`$-verletzenden Prozessen, welche bei $`T0`$ durch einen Faktor $`\mathrm{exp}(2\pi /\alpha _2)`$ unterdrückt sind und vernachlässigt werden können. Bei Temperaturen $`10^2`$ GeV $`T`$ $`10^{12}`$ GeV jedoch sind sie in Gegenwart von statischen topologischen Feldkonfigurationen, den sogenannten Sphaleronen, stark genug, um eine vorhandene $`B+L`$-Asymmetrie vollständig auszuwaschen . Andererseits sind sie in der Lage, eine bereits vorhandene $`BL`$-Asymmetrie teilweise in eine $`B`$-Asymmetrie gemäß
$$B(TM_Z)=\frac{32+4n_H}{98+13n_H}(BL)$$
(4.89)
umzuwandeln.
Nun erhalten im Rahmen der $`BL`$-Symmetriebrechung bei $`M_I`$ die schweren Neutrinos ihre Massen. Aufgrund ihrer Majorana-Natur können sie über leptonzahlverletzende Prozesse
$$N_Rl_L+\overline{\varphi },N_Rl_L^C+\varphi $$
(4.90)
sowohl in Leptonen als auch in Antileptonen zerfallen; $`\varphi `$ bezeichnet ein Higgs-Teilchen. Die notwendige $`CP`$-Verletzung wird durch komplexe Kopplungen zwischen $`N_R`$, $`l_L`$ und $`\overline{\varphi }`$ realisiert. Auf diese Weise erhält man eine reine Lepton- und somit eine $`BL`$-Asymmetrie, die gemäß (4.89) in die beobachtete Baryonasymmetrie des Universums konvertiert werden kann . Damit die leptonzahlverletzenden Zerfälle der schweren Neutrinos im erforderlichen thermischen Ungleichgewicht stattfinden, muß die zugehörige Zerfallsrate kleiner als die Expansionsrate des Universums bei dieser Temperatur sein. Das läßt sich qualitativ als Bedingung
$$m_\nu \left(\frac{M}{10^{10}\text{GeV}}\right)^{1/2}4\text{eV}$$
(4.91)
für die Massen $`m_\nu `$ der leichten Neutrinos schreiben, wobei $`M`$ Min$`\{M_I\mathrm{,10}^{12}`$ GeV$`\}`$ ist . Im hier untersuchten Massenmodell ist (4.91) offensichtlich erfüllt.
## Kapitel 5 Zerfallsraten der Nukleonen
Wie in Kapitel 2 dargestellt wurde, gehört die Instabilität der Nukleonen zu den bemerkenswertesten Vorhersagen von GUTs. Die Bestimmung der partiellen und totalen Zerfallsraten im Rahmen des hier untersuchten $`SO(10)`$-Massenmodells wird Gegenstand dieses Kapitels sein. Dazu sind zwei wesentliche Schritte notwendig:
* Zunächst muß die vollständige effektive Lagrangedichte $`_{\text{eff}}`$ der durch Eichbosonaustausch vermittelten Nukleonenzerfälle in der Basis der Masseneigenzustände bestimmt werden. Das geschieht analog zur Herleitung der Lagrangedichte der Fermi-Theorie aus jener der elektroschwachen Wechselwirkung.
* Weiterhin muß man die hadronischen Übergangsmatrixelemente berechnen. Hierzu sind in der Vergangenheit das nichtrelativistische Quark-Modell , das Bag-Modell und die chirale Störungsrechnung verwendet worden. Ursprüngliche Diskrepanzen zwischen den Resultaten der chiralen Störungsrechnung und denen der Quark-Modelle sind in beseitigt worden; gibt eine Übersicht über die mit den verschiedenen Verfahren erzielten Ergebnisse. In dieser Arbeit wird analog zu und vorgegangen.
Prinzipiell existieren auch Nukleonzerfallsprozesse aufgrund von Higgs-Austausch. Diese sind aber selbst dann vernachlässigbar, wenn das Higgs-Teilchen eine Masse der Größenordnung $`M_I`$ hat, da die Übergangswahrscheinlichkeit proportional zur vierten Potenz der sehr kleinen Yukawa-Kopplungen an die drei leichten Quarks ist.
### 5.1 Die effektive Lagrangedichte
Ausgangspunkt ist der baryonzahlverletzende Anteil der $`SO(10)`$-Lagrangedichte (2.27). Verallgemeinert man ihn auf drei Fermiongenerationen und ersetzt die Wechselwirkungs- durch die Masseneigenzustände analog zu (1.30), kann man die vollständige effektive Lagrangedichte der Vier-Fermion-Wechselwirkungen berechnen, welche die eichbosonvermittelten Nukleonenzerfälle beschreibt. Für den Fall zweier Generationen ist das in durchgeführt worden. Unter Benutzung der Fierz-Identitäten
$`\left(\overline{\mathrm{\Psi }}_{1L}\gamma _\mu \mathrm{\Psi }_{2L}\right)\left(\overline{\mathrm{\Psi }}_{3L}\gamma ^\mu \mathrm{\Psi }_{4L}\right)`$ $`=`$ $`\left(\overline{\mathrm{\Psi }}_{1L}\gamma _\mu \mathrm{\Psi }_{4L}\right)\left(\overline{\mathrm{\Psi }}_{3L}\gamma ^\mu \mathrm{\Psi }_{2L}\right),`$ (5.1)
$`\overline{\mathrm{\Psi }}_{1L}\gamma _\mu \mathrm{\Psi }_{2L}`$ $`=`$ $`\overline{\mathrm{\Psi }}_{2R}^C\gamma _\mu \mathrm{\Psi }_{1R}^C`$ (5.2)
erhält man als Resultat:
$`_{\text{eff}}`$ $`=`$ $`A_1\left(\epsilon _{\alpha \beta \gamma }\overline{u}_L^{C\gamma }\gamma ^\mu u_L^\beta \right)\left(\overline{e}_L^+\gamma _\mu d_L^\alpha \right)+A_2\left(\epsilon _{\alpha \beta \gamma }\overline{u}_L^{C\gamma }\gamma ^\mu u_L^\beta \right)\left(\overline{e}_R^+\gamma _\mu d_R^\alpha \right)`$ (5.3)
$`+`$ $`A_3\left(\epsilon _{\alpha \beta \gamma }\overline{u}_L^{C\gamma }\gamma ^\mu u_L^\beta \right)\left(\overline{\mu }_L^+\gamma _\mu d_L^\alpha \right)+A_4\left(\epsilon _{\alpha \beta \gamma }\overline{u}_L^{C\gamma }\gamma ^\mu u_L^\beta \right)\left(\overline{\mu }_R^+\gamma _\mu d_R^\alpha \right)`$
$`+`$ $`A_5\left(\epsilon _{\alpha \beta \gamma }\overline{u}_L^{C\gamma }\gamma ^\mu u_L^\beta \right)\left(\overline{e}_L^+\gamma _\mu s_L^\alpha \right)+A_6\left(\epsilon _{\alpha \beta \gamma }\overline{u}_L^{C\gamma }\gamma ^\mu u_L^\beta \right)\left(\overline{e}_R^+\gamma _\mu s_R^\alpha \right)`$
$`+`$ $`A_7\left(\epsilon _{\alpha \beta \gamma }\overline{u}_L^{C\gamma }\gamma ^\mu u_L^\beta \right)\left(\overline{\mu }_L^+\gamma _\mu s_L^\alpha \right)+A_8\left(\epsilon _{\alpha \beta \gamma }\overline{u}_L^{C\gamma }\gamma ^\mu u_L^\beta \right)\left(\overline{\mu }_R^+\gamma _\mu s_R^\alpha \right)`$
$`+`$ $`A_9\left(\epsilon _{\alpha \beta \gamma }\overline{u}_L^{C\gamma }\gamma ^\mu d_L^\beta \right)\left(\overline{\nu }_{eR}^C\gamma _\mu d_R^\alpha \right)+A_{10}\left(\epsilon _{\alpha \beta \gamma }\overline{u}_L^{C\gamma }\gamma ^\mu d_L^\beta \right)\left(\overline{\nu }_{\mu R}^C\gamma _\mu d_R^\alpha \right)`$
$`+`$ $`A_{11}\left(\epsilon _{\alpha \beta \gamma }\overline{u}_L^{C\gamma }\gamma ^\mu d_L^\beta \right)\left(\overline{\nu }_{eR}^C\gamma _\mu s_R^\alpha \right)+A_{12}\left(\epsilon _{\alpha \beta \gamma }\overline{u}_L^{C\gamma }\gamma ^\mu d_L^\beta \right)\left(\overline{\nu }_{\mu R}^C\gamma _\mu s_R^\alpha \right)`$
$`+`$ $`A_{13}\left(\epsilon _{\alpha \beta \gamma }\overline{u}_L^{C\gamma }\gamma ^\mu s_L^\beta \right)\left(\overline{\nu }_{eR}^C\gamma _\mu d_R^\alpha \right)+A_{14}\left(\epsilon _{\alpha \beta \gamma }\overline{u}_L^{C\gamma }\gamma ^\mu s_L^\beta \right)\left(\overline{\nu }_{\mu R}^C\gamma _\mu d_R^\alpha \right)`$
$`+`$ $`A_{15}\left(\epsilon _{\alpha \beta \gamma }\overline{u}_L^{C\gamma }\gamma ^\mu d_L^\beta \right)\left(\overline{\nu }_{\tau R}^C\gamma _\mu d_R^\alpha \right)+A_{16}\left(\epsilon _{\alpha \beta \gamma }\overline{u}_L^{C\gamma }\gamma ^\mu d_L^\beta \right)\left(\overline{\nu }_{\tau R}^C\gamma _\mu s_R^\alpha \right)`$
$`+`$ $`A_{17}\left(\epsilon _{\alpha \beta \gamma }\overline{u}_L^{C\gamma }\gamma ^\mu s_L^\beta \right)\left(\overline{\nu }_{\tau R}^C\gamma _\mu d_R^\alpha \right)`$
$`+`$ ( Terme mit zwei $`s`$-Quarks )
$`+`$ ( Terme mit $`c`$\- ,$`b`$\- und $`t`$-Quarks )
$`+`$ ( Terme mit $`\overline{\tau }_{L,R}^+`$ und $`\overline{\nu }_{e,\mu ,\tau L}^C`$ )
$`+`$ ( h.c. )
Die $`A_i`$-Koeffizienten sind in Anhang E.1 aufgeführt und hängen außer von den Eichbosonmassen $`M_{X,Y}`$ und $`M_{X^{},Y^{}}`$, für welche $`M_{X,Y}=M_{X^{},Y^{}}M_U`$ angenommen wird, von $`g_U`$ und den Elementen der fermionischen Mischungsmatrizen $`𝐔_{R,L}`$, $`𝐃_{R,L}`$, $`𝐄_{R,L}`$ und $`𝐍_L`$ ab. Diese Größen sind bereits im vorigen Kapitel bestimmt worden, so daß alle $`A_i`$ bekannt sind. $`𝐍_R`$ hat keinen Einfluß auf die $`A_i`$, da sie die Massenmatrix der schweren Neutrinos diagonalisiert, welche als Zerfallsprodukte der Nukleonen nicht vorkommen. An dieser Stelle wird klar, warum die Fermionmischungen für die Verzweigungsraten der Nukleonenzerfälle von großer Bedeutung sind.
### 5.2 Hadronische Übergangsmatrixelemente
Die hier verwendete Methode zur Berechnung der Übergangsmatrixelemente entspricht dem sogenannten „recoil“-Modell in . Grundlage ist ein Quark-Modell mit $`SU(6)`$-Spin-Flavour-Symmetrie; die Quarks treten in drei Flavours $`u`$, $`d`$ und $`s`$ sowie zwei Spinzuständen $``$ und $``$ auf. Alle sechs unabhängigen Quarkzustände haben dieselbe Masse $`m`$, was eine einheitliche Masse für die Mesonen zur Folge hat. Die $`SU(6)`$-Symmetrie wird später bei der Berechnung der Zerfallsraten durch einen Phasenraumfaktor explizit gebrochen.
Es wird in einem Bezugssystem gearbeitet, in dem das Nukleon, bevor es zerfällt, im Ursprung des Systems in Ruhe ist. Die Spineinstellungen $`+1/2`$ beziehungsweise $`1/2`$ werden relativ zur $`z`$-Achse angegeben. Der Zerfall führt stets auf ein Antilepton, welches sich in positive $`z`$-Richtung bewegt, während das Meson einen Rückstoß in Richtung der negativen $`z`$-Achse erhält. Das entstehende Antilepton wird als relativistisch angenommen und es gilt $`l_R^Cl^C`$ und $`l_L^Cl^C`$.
Die Vorgehensweise zur Berechnung der Matrixelemente soll kurz am Beispiel von $`e^+\pi ^0|_{\text{eff}}|p`$ erläutert werden, für Details sei auf verwiesen. Aufgrund der Drehimpulserhaltung läßt sich die Übergangswahrscheinlichkeit zunächst gemäß
$$|e^+\pi ^0|_{\text{eff}}|p|^2|e_R^+\pi ^0|_{\text{eff}}|p|^2+|e_L^+\pi ^0|_{\text{eff}}|p|^2$$
(5.4)
zerlegen. Eine weitere Zerlegung ist möglich, indem der Zustand $`|\pi ^0`$ durch den Spin-, Flavour- und Farbanteil seiner Wellenfunktion ausgedrückt wird:
$$|\pi ^0=\mathrm{\hspace{0.33em}1}/\sqrt{12}\delta _{\rho \sigma }|u^{C\rho }u^\sigma u^{C\rho }u^\sigma d^{C\rho }d^\sigma +d^{C\rho }d^\sigma $$
(5.5)
Dadurch ergeben sich die Matrixelemente der für den Zerfall $`pe^+\pi ^0`$ relevanten Elementarprozesse:
$`e_R^+u^Cu|_{\text{eff}}|p`$ , $`e_R^+u^Cu|_{\text{eff}}|p,`$
$`e_R^+d^Cd|_{\text{eff}}|p`$ , $`e_R^+d^Cd|_{\text{eff}}|p`$ (5.6)
Die Matrixelemente mit $`e_L^+`$ im Endzustand erhält man analog. Auch das Proton kann gemäß
$$|p=\mathrm{\hspace{0.33em}1}/\sqrt{18}\epsilon _{\alpha \beta \gamma }|(u^\alpha d^\beta u^\alpha d^\beta )u^\gamma $$
(5.7)
in seine Spin-, Flavour- und Farbanteile zerlegt werden. Die Ortsanteile der Wellenfunktionen werden später berücksichtigt, da sie in guter Näherung für alle Zerfallsprozesse gleich groß sind. Die Spin-Flavour-Anteile der Mesonwellenfunktionen können in Anhang E.2 gefunden werden.
Weiterhin identifiziert man die für die Bestimmung von (5.4) relevanten Terme in $`_{\text{eff}}`$ als
* $`A_2\left(\epsilon _{\alpha \beta \gamma }\overline{u}_L^{C\gamma }\gamma ^\mu u_L^\beta \right)\left(\overline{e}_R^+\gamma _\mu d_R^\alpha \right)`$ für $`e_R^+\pi ^0|_{\text{eff}}|p`$
* $`A_1\left(\epsilon _{\alpha \beta \gamma }\overline{u}_L^{C\gamma }\gamma ^\mu u_L^\beta \right)\left(\overline{e}_L^+\gamma _\mu d_L^\alpha \right)`$ für $`e_L^+\pi ^0|_{\text{eff}}|p`$
Im nächsten Schritt werden nun alle Quark- und Antiquarkfelder sowohl in den Zustandsfunktionen $`|p`$ und $`\pi ^0|`$ als auch in den Lagrangedichtetermen in Form von Erzeugungs- und Vernichtungsoperatoren sowie Dirac-Spinoren u und v ausgedrückt. Für ein Quark $`q`$ beziehungsweise Antiquark $`q^C`$ erhält man auf diese Weise
$$q\widehat{a}(q)𝐮+\widehat{b}^{}(q)𝐯,q^C\widehat{a}(q)𝐯+\widehat{b}^{}(q)𝐮,$$
(5.8)
wobei der Operator $`\widehat{a}(q)`$ ein $`q`$-Quark vernichtet und der Operator $`\widehat{b}^{}(q)`$ ein $`q^C`$-Antiquark erzeugt; die Spinoren u und v mit
$$(p/m)𝐮=\mathrm{\hspace{0.33em}0},(p/+m)𝐯=\mathrm{\hspace{0.33em}0}$$
(5.9)
sind in definiert. Man erhält auf diese Weise umfangreiche Ausdrücke für die Größen in (5.6), welche sich jedoch beträchtlich vereinfachen, wenn man unter Ausnutzung der Antikommutator-Eigenschaften
$`\{\widehat{a}(q),\widehat{a}^{}(\stackrel{~}{q})\}=\{\widehat{b}(q),\widehat{b}^{}(\stackrel{~}{q})\}\delta _{q\stackrel{~}{q}}`$ (5.10)
$`\{\widehat{a}(q),\widehat{a}(\stackrel{~}{q})\}=\{\widehat{b}(q),\widehat{b}(\stackrel{~}{q})\}=\{\widehat{a}^{}(q),\widehat{a}^{}(\stackrel{~}{q})\}=\{\widehat{b}^{}(q),\widehat{b}^{}(\stackrel{~}{q})\}=\mathrm{\hspace{0.33em}0}`$ (5.11)
$`\{\widehat{a}(q),\widehat{b}(\stackrel{~}{q})\}=\{\widehat{a}^{}(q),\widehat{b}(\stackrel{~}{q})\}=\{\widehat{a}(q),\widehat{b}^{}(\stackrel{~}{q})\}=\{\widehat{a}^{}(q),\widehat{b}^{}(\stackrel{~}{q})\}=\mathrm{\hspace{0.33em}0}`$ (5.12)
die Vernichtungsoperatoren auf den Vakuumzustand $`|0`$ und die Erzeugungsoperatoren auf den Zustand $`0|`$ wirken läßt und dadurch für die meisten Terme Null erhält. In den nichtverschwindenden Beiträgen verbleiben die Spinor-Amplituden, welche noch berechnet werden müssen.
Im „recoil“-Modell wird vom Übergang zweier statischer Quarks derselben Masse $`m`$ in ein relativistisches Antilepton und ein Antiquark, welches einen Rückstoßimpuls $`p`$ mit $`(p/m=3/4)`$ besitzt, ausgegangen. Die Spinoren der Antileptonen sind demnach durch
$`l_R^C={\displaystyle \frac{1}{\sqrt{2}}}\left(\begin{array}{c}\chi \\ \chi \end{array}\right)`$ , $`l_L^C={\displaystyle \frac{1}{\sqrt{2}}}\left(\begin{array}{c}\chi \\ \chi \end{array}\right)`$ (5.13)
und die der statischen Quarks durch
$$𝐮=\left(\begin{array}{c}\chi \\ 0\end{array}\right),𝐮=\left(\begin{array}{c}\chi \\ 0\end{array}\right),𝐯=\left(\begin{array}{c}0\\ \chi \end{array}\right),𝐯=\left(\begin{array}{c}0\\ \chi \end{array}\right)$$
(5.14)
gegeben. Die nichtstatischen Antiquarks im Endzustand werden durch die Spinoren
$`𝐮={\displaystyle \frac{1}{\sqrt{10}}}\left(\begin{array}{c}3\chi \\ \chi \end{array}\right)`$ , $`𝐮={\displaystyle \frac{1}{\sqrt{10}}}\left(\begin{array}{c}3\chi \\ \chi \end{array}\right),`$
$`𝐯={\displaystyle \frac{1}{\sqrt{10}}}\left(\begin{array}{c}\chi \\ 3\chi \end{array}\right)`$ , $`𝐯={\displaystyle \frac{1}{\sqrt{10}}}\left(\begin{array}{c}\chi \\ 3\chi \end{array}\right)`$ (5.15)
beschrieben. Nun lassen sich die Spinor-Amplituden bestimmen; Tabelle 5.1 enthält alle nichtverschwindenden Amplituden und deren Werte.
Die Amplituden zu den Zerfällen mit linkshändigen Antileptonen im Endzustand erhält man durch Ersetzen von $`\overline{l}_R^C`$ mit $`\overline{l}_L^C`$ und Vertauschen von $``$ und $``$.
Damit kann man die Matrixelemente (5.6) berechnen, und aus diesen lassen sich unter Berücksichtigung der Spin-Flavour-Struktur (5.5) des Mesons wiederum die Übergangswahrscheinlichkeiten in (5.4) bestimmen. Analog verfährt man für alle relevanten Zerfallskanäle; die Tabelle E.4 gibt die Resultate für die Amplituden der elementaren Zerfallsprozesse an, während Tabelle E.5 die Übergangswahrscheinlichkeiten für die Zerfälle mit physikalischen Endzuständen beinhaltet. Zu beachten ist die Tatsache, daß die Amplituden der einzelnen Elementarprozesse von unterschiedlichen $`A_i`$-Koeffizienten gemäß (5.3) begleitet werden.
### 5.3 Berechnung der Zerfallsraten
Für die partielle Zerfallsrate eines bestimmten Prozesses $`\text{Nukleon}\text{Meson}+\text{Antilepton}`$ gilt die Beziehung
$$\mathrm{\Gamma }_j=\frac{1}{16\pi }m_{\text{Nukl}}^2\rho _j|S|^2|𝒜|^2\left(|𝒜_L|^2\underset{l}{}|A_l_l|^2+|𝒜_R|^2\underset{r}{}|A_r_r|^2\right),$$
(5.16)
welche man aus der allgemeinen Bestimmungsgleichung für $`\mathrm{\Gamma }`$ in Zwei-Teilchen-Zerfällen unter Berücksichtigung der spezifischen Annahmen des „recoil“-Modells erhält. (5.16) ist äquivalent zu der in gegebenen Gleichung für $`\mathrm{\Gamma }_j`$ und unterscheidet sich von dieser lediglich darin, daß der Phasenraumfaktor und die Übergangsmatrixelemente aus verwendet wurden, welche direkt proportional zu den in benutzten sind. Weiterhin ist in (5.16) die Abhängigkeit von $`\alpha _U`$ und den Eichbosonmassen in den Koeffizienten $`A_i`$ aus (5.3) enthalten. Die verschiedenen Größen in (5.16) sind wie folgt definiert:
* $`_l`$ und $`_r`$ sind die für den jeweiligen Zerfallsprozeß relevanten hadronischen Übergangsmatrixelemente, wobei sich die Indizes $`l`$ und $`r`$ auf die Chiralität des Antileptons im Endzustand beziehen. Die Summation berücksichtigt die Tatsache, daß für einige Zerfälle zwei Lagrangedichteterme verantwortlich sind. Die Bestimmung der Matrixelemente war Gegenstand des letzten Abschnitts.
* $`A_l`$ und $`A_r`$ sind die zu $`_l`$ beziehungsweise $`_r`$ gehörenden Koeffizienten aus der effektiven Lagrangedichte. Gemäß (5.3) treten alle Matrixelemente der elementaren Zerfallsprozesse in Kombination mit einem bestimmten $`A_i`$-Koeffizient auf. Tabelle E.5 gibt die Werte für die unabhängigen $`|A_r_r|^2`$ an; alle dort nicht explizit aufgeführten Größen lassen sich durch Symmetrieüberlegungen aus diesen herleiten.
* $`𝒜,𝒜_L`$ und $`𝒜_R`$ sind Faktoren, welche aus der Renormierung der Vier-Fermion-Operatoren in der effektiven Lagrangedichte resultieren. Während nämlich die Berechnung der hadronischen Matrixelemente $`_{r,l}`$ für Energieskalen $`\mu m_p`$ durchgeführt wird, ist die eigentliche baryonzahlverletzende Wechselwirkung bei Skalen $`\mu M_U`$ wirksam. Das führt zu Renormierungseffekten, welche auf Strahlungskorrekturen zu den Vier-Fermion-Operatoren durch die $`G_{\text{PS}}`$\- und SM-Eichbosonen beruhen. Eine formale Analyse dieses Phänomens im Rahmen einer $`SU(5)`$-GUT ist in zu finden, während die Renormierungseffekte für den Fall einer $`SO(10)`$-Theorie untersucht. Für das in dieser Arbeit untersuchte Modell sind die Faktoren durch
$`𝒜_L`$ $`=`$ $`\left({\displaystyle \frac{\alpha _1(M_Z)}{\alpha _1(M_I)}}\right)^{\frac{23}{82}}`$ (5.17)
$`𝒜_R`$ $`=`$ $`\left({\displaystyle \frac{\alpha _1(M_Z)}{\alpha _1(M_I)}}\right)^{\frac{11}{82}}`$ (5.18)
$`𝒜`$ $`=`$ $`\left({\displaystyle \frac{\alpha _{4C}(M_I)}{\alpha _{4C}(M_U)}}\right)^{\frac{5}{8}}\left({\displaystyle \frac{\alpha _{2L}(M_I)}{\alpha _{2L}(M_U)}}\right)^{\frac{27}{100}}\left({\displaystyle \frac{\alpha _{2R}(M_I)}{\alpha _{2R}(M_U)}}\right)^{\frac{3}{20}}\left({\displaystyle \frac{\alpha _2(M_Z)}{\alpha _2(M_I)}}\right)^{\frac{27}{38}}`$ (5.19)
$`\left({\displaystyle \frac{\alpha _3(M_Z)}{\alpha _3(M_I)}}\right)^{\frac{2}{7}}\left({\displaystyle \frac{\alpha _3(m_b)}{\alpha _3(M_Z)}}\right)^{\frac{6}{23}}\left({\displaystyle \frac{\alpha _3(m_c)}{\alpha _3(m_b)}}\right)^{\frac{6}{25}}\left({\displaystyle \frac{\alpha _3(1\text{GeV})}{\alpha _3(m_c)}}\right)^{\frac{2}{9}}`$
gegeben. In die Exponenten gehen die anomalen Dimensionen der Vier-Fermion-Operatoren und die führenden Koeffizienten der jeweiligen Eichkopplungs-$`\beta `$-Funktionen ein. Die Eichkopplungen bei $`M_Z`$ sind in Tabelle C.1 zu finden, diejenigen bei $`M_U`$ und $`M_I`$ wurden im vorigen Kapitel bestimmt; Anhang C.2 enthält die numerischen Resultate. Die Werte von $`\alpha _3`$ bei Energien $`\mu <M_Z`$ sind entnommen:
$$\alpha _3(1\text{GeV})=\mathrm{\hspace{0.33em}0.544},\alpha _3(m_c)=\mathrm{\hspace{0.33em}0.412},\alpha _3(m_b)=\mathrm{\hspace{0.33em}0.226}$$
(5.20)
Der Einfluß der Skalenabhängigkeit von $`\alpha _{\text{em}}`$ wird vernachlässigt, da er viel kleiner als der von $`\alpha _3`$ ist ($`\alpha _{\text{em}}^1(M_Z)129`$, $`\alpha _{\text{em}}^1(0)137`$). Unter Verwendung aller Kopplungswerte erhält man
$$|𝒜_L|^2=\mathrm{\hspace{0.33em}1.155},|𝒜_R|^2=\mathrm{\hspace{0.33em}1.071},|𝒜|^2=\mathrm{\hspace{0.33em}23.59}$$
(5.21)
* $`|S|^2=\mathrm{\Psi }_{\text{Nucl}}^s(\stackrel{}{r}_1,\stackrel{}{r}_2,\stackrel{}{r}_3)|\delta (\stackrel{}{r}_1\stackrel{}{r}_2)|\mathrm{\Psi }_{\text{Nucl}}^s(\stackrel{}{r}_1,\stackrel{}{r}_2,\stackrel{}{r}_3)`$ steht für die Wahrscheinlichkeit, zwei Valenzquarks des Nukleons in einem Raumpunkt zu finden. Das ist erforderlich, damit der Elementarprozeß $`qqq^cl^c`$ stattfinden kann, da die zugrundeliegende Wechselwirkung extrem kurzreichweitig ist. $`\mathrm{\Psi }_{\text{Nucl}}^s`$ ist der Raumanteil der Wellenfunktion des Nukleons. Der räumliche Anteil der Mesonwellenfunktion wird wie in vernachlässigt; die Spin-, Flavour- und Farbanteile der Wellenfunktionen sind bereits in die Berechnung von $`_{r,l}`$ eingegangen. Für $`|S|^2`$ wird der Wert 0.012 GeV<sup>3</sup> verwendet .
* $`\rho _j(1\chi _j^2)(1\chi _j^4)`$ mit $`\chi _j=m_{\text{Meson}}/m_{\text{Nukl}}`$ ist ein Phasenraumfaktor, welcher die $`SU(6)`$-Spin-Flavour-Symmetrie des im letzten Abschnitt verwendeten Quarkmodells explizit bricht. Er berücksichtigt die in der Natur existierenden unterschiedlichen Massen der Mesonen in den Endzuständen der Zerfallsprozesse und ihren Einfluß auf die Raten. Tabelle E.3 enthält die relevanten Werte von $`\rho _p`$ und $`\rho _n`$.
Damit sind alle Größen in (5.16) bekannt und die partiellen sowie totalen Zerfallsraten der Nukleonen können berechnet werden. Sämtliche Resultate sowohl für die untersuchten Modelle als auch für den Fall verschwindender Fermionmischungen sind in den Tabellen E.7-E.10 im Anhang zusammengefaßt.
Man erkennt die generelle Tendenz, daß der Einfluß der Mischungen zu einer deutlichen Unterdrückung des Zerfallskanals $`p,ne^+X`$ führt, während die Kanäle $`p,n\mu ^+X`$ und $`p,n\nu ^CX`$ im Vergleich zum Fall verschwindender Mischungen bevorzugt werden. Besonders auffällig ist das bei den Zerfällen $`p,n\mu ^+\pi `$, $`p,n\nu ^CK`$ und $`p\mu ^+\omega `$. Insofern unterscheiden sich die Verzweigungsraten zum Teil beträchtlich von denen nichtsupersymmetrischer GUT-Modelle mit kleinen Mischungen, welche eine deutliche Dominanz von $`p,ne^+X`$-Zerfällen vorhersagen. SUSY-GUTs dagegen bevorzugen Zerfälle mit $`K`$-Mesonen im Endzustand .
Zwischen den Verzweigungsraten der drei analysierten Modelle bestehen teilweise deutlich ausgeprägte Unterschiede, was man auch an den in Tabelle 5.2 dargestellten Verhältnissen partieller Raten erkennen kann.
Während die Modelle 2a und 2b anhand ihrer Eigenschaften im Neutrinosektor praktisch nicht unterscheidbar waren, zeigen sich in ihren Verzweigungsraten Differenzen, welche durchaus experimentell zugänglich sein sollten, sofern in Zukunft Nukleonenzerfälle beobachtet werden. Die totalen Zerfallsraten der einzelnen Modelle sind nahezu gleich und führen zu Lebensdauern im Bereich von $`\tau _{p,n}(34)10^{34}`$ Jahren.
Die in den Tabellen E.7-E.10 dargelegten Resultate dieser Untersuchung sind echte Vorhersagen des verwendeten $`SO(10)`$-Massenmodells und verdeutlichen auf eindrucksvolle Weise den Einfluß der Fermionmischungen auf die Verzweigungsraten der Nukleonenzerfälle.
Bei allen Rechnungen ist stets mit den Mittelwerten der verwendeten Größen gearbeitet worden, ohne deren Fehler explizit zu berücksichtigen. Das ist insofern begründet, als die Hauptquellen der auftretenden Unsicherheiten in den Einflüssen der Schwellenkorrekturen und den modellspezifischen Näherungen bei der Bestimmung der hadronischen Matrixelemente zu sehen sind. Gerade diese sind aber im Gegensatz zu den bekannten Meßungenauigkeiten in den SM-Kopplungen und den Fermionmassen und -mischungen nur schwer abzuschätzen. Für die Schwelleneffekte und deren Einluß auf die Werte der Symmetriebrechungsskalen ist das in Abschnitt 4.2 versucht worden, während die Auswirkung der Unsicherheit in den Übergangsmatrixelementen auf die Lebensdauer des Protons mit einem Faktor $`10^{\pm 0.7}`$ ansetzt. Es kann allerdings davon ausgegangen werden, daß die Effekte dieser Unsicherheiten in erster Linie die totalen Zerfallsraten beeinflussen, die Verzweigungsraten dagegen deutlich weniger betreffen. Letztere stellen aber die wirklich relevanten Vorhersagen des Massenmodells dar.
### 5.4 Experimenteller Status
Es bleibt zu prüfen, ob die im letzten Abschnitt gemachten Vorhersagen bezüglich der partiellen und totalen Zerfallsraten der Nukleonen mit den experimentellen Grenzen für diese Größen verträglich sind.
Die Experimente, welche nach Nukleonenzerfällen suchen, lassen sich in zwei Kategorien einteilen. Wasser-Čerenkovdetektoren wie IMB 3 und Kamiokande beziehungsweise Super-Kamiokande sind besonders zum Nachweis von $`p,ne^+\pi `$ geeignet, während die Eisenkalorimeter-Spurdetektoren wie Fréjus und Soudan 2 für Zerfälle mit $`K`$-Mesonen und $`\mu ^+`$ im Endzustand sensitiv sind. Keines der Experimente hat bis heute eindeutige Hinweise auf die Instabilität der Nukleonen liefern können, so daß lediglich obere Grenzen für die Zerfallsraten existieren.
Die in angegebenen Werte für diese Grenzen sind in Tabelle E.6 aufgeführt, seitdem veröffentlichte Resultate von Super-Kamiokande zeigt Tabelle 5.3. Der Vergleich mit den Modellvorhersagen in den Tabellen E.7-E.10 führt zu dem Schluß, daß diese mit allen experimentellen Grenzen verträglich sind.
Die Grenzwerte für die Lebensdauern bezüglich bestimmter Zerfallsprozesse wie zum Beispiel $`pe^+\pi ^0,\mu ^+\pi ^0`$ erreichen gemäß Tabelle 5.3 bereits den Bereich einiger $`10^{33}`$ Jahre, was lediglich eine Größenordnung unter den in dieser Arbeit gemachten Vorhersagen liegt. Letztere werden also in absehbarer Zeit direkt überprüfbar sein, da Super-Kamiokande weiterhin Daten aufnimmt und Experimente wie das auf einem Detektor aus flüssigem Argon basierende ICARUS neu hinzukommen werden.
## Zusammenfassung und Ausblick
Gegenstand dieser Untersuchung war ein Modell für die fermionischen Massenmatrizen auf der Grundlage einer $`SO(10)`$-Theorie, welche über eine intermediäre Pati-Salam-Symmetrie in das Standardmodell gebrochen wird. Der gewählte Ansatz bestand in einer asymmetrischen „Nearest Neighbour Interaction“-Form der Matrizen, die durch eine globale $`U(1)`$-Familiensymmetrie realisiert wird. Dieser beinhaltet im Rahmen des Standardmodells keine physikalischen Konsequenzen, da die rechtshändigen Mischungen dort nicht beobachtbar sind. In Theorien jenseits des Standardmodells jedoch haben alle Fermionmischungen beträchtlichen Einfluß auf observable Größen wie die Zerfallsraten der Nukleonen.
Nach der systematischen Bestimmung der Symmetriebrechungsskalen und Eichkopplungen in Abhängigkeit vom Teilcheninhalt des Modells sowie einer Abschätzung der Schwelleneffekte wurden die Massenmatrizen unter Berücksichtigung der $`SO(10)`$-spezifischen Beziehungen zwischen diesen konstruiert. Anschließend sind alle Lösungen des Modells numerisch bestimmt worden, welche neben den bekannten Parametern im Sektor der geladenen Fermionen auch phänomenologisch sinnvolle Neutrinoeigenschaften liefern. Es wurden insgesamt drei Lösungen des Modells gefunden, welche in der Lage sind, die Anomalien der Sonnen- und atmosphärischen Neutrinos durch Oszillationen zu erklären. Zwei davon deuten auf eine Lösung des Sonnen-Neutrinoproblems durch den MSW-Effekt mit kleiner Mischung hin, während die dritte den MSW-Effekt mit großer Mischung beinhaltet. Bemerkenswert ist, daß alle Lösungen auf leptonische (23)-Mischungen führen, die im Rahmen des zur Erklärung der Anomalie atmosphärischer Neutrinos erlaubten Parameterbereichs vergleichsweise klein sind. Da die (23)-Mischung desjenigen Modells, welches den MSW-Effekt mit großer Mischung realisiert, nahe an der unteren Grenze des zulässigen Bereichs liegt und letzterer durch die Oszillationsexperimente in naher Zukunft weiter eingeschränkt werden wird, sind die anderen beiden Lösungen als realistischer zu betrachten. Im Rahmen des hier analysierten Massenmodells erscheint eine Lösung des Sonnen-Neutrinoproblems durch den MSW-Effekt mit kleiner Mischung als wahrscheinlich.
Alle gefundenen Lösungen besitzen mehrere betragsmäßig große Mischungen im Sektor der geladenen Fermionen. Dabei handelt es sich sowohl um im Standardmodell nicht observable rechtshändige als auch um linkshändige Mischungswinkel. Demnach kann man von der Größe der CKM-Winkel nicht zwangsläufig auf die der linkshändigen Fermionmischungen schließen, wie es häufig getan wird. Die teilweise großen Mischungen führen zu charakteristischen Verzweigungsraten der Nukleonenzerfälle, welche von denen bei verschwindenden Mischungen deutlich abweichen. Generell läßt sich eine Unterdrückung der Kanäle mit $`e^+`$ im Endzustand feststellen, während Kanäle mit $`\mu ^+`$ und $`\nu ^C`$ bevorzugt werden. Die drei Lösungen sind aufgrund ihrer Vorhersagen für die Verzweigungsraten auch untereinander unterscheidbar. Die berechneten Lebensdauern der Nukleonen betragen $`\tau _{p,n}(34)10^{34}`$ Jahre, was je nach Zerfallskanal lediglich ein bis zwei Größenordnungen über den experimentellen Grenzen liegt. Damit sind die Modellvorhersagen in absehbarer Zeit direkt überprüfbar.
Der Einfachheit halber sind alle Massenmatrizen als reell angenommen worden, da das Problem der $`CP`$-Verletzung nicht Gegenstand der Arbeit war. Abgesehen von der Standardmethode, die $`CP`$-Verletzung über komplexe Yukawa-Kopplungen zu realisieren, bieten Grand Unified-Theorien aufgrund ihres erweiterten Eichboson- und Higgs-Sektors dafür verschiedene weitere Möglichkeiten .
Weiterhin wurde auf die Verwendung von Supersymmetrie verzichtet, da supersymmetrische Theorien wegen des wesentlich umfangreicheren Teilchenspektrums zahlreiche zusätzliche Parameter enthalten, deren Werte weitgehend unbestimmt sind. Das hat naturgemäß negative Auswirkungen auf die Vorhersagekraft der Modelle. Auch eine für die Anwendbarkeit des See-Saw-Mechanismus vorteilhafte intermediäre Skala im Bereich $`10^{10}10^{12}`$ GeV ist in supersymmetrischen Theorien nicht auf natürliche Weise realisierbar, da sich die Eichkopplungen des supersymmetrischen Standardmodells im Gegensatz zum nichtsupersymmetrischen Fall genau in einem Punkt treffen. Die Abwesenheit jeglicher experimenteller Hinweise auf eine bei Energien im TeV-Bereich gebrochene Supersymmetrie in der Natur ist ein weiterer Schwachpunkt, mit dem supersymmetrische Theorien bisher behaftet sind. Schließlich deutet auch die Existenz nichtsupersymmetrischer Niederenergie-Vakua in der M-Theorie darauf hin, daß Supersymmetrie keine zwingende Eigenschaft von Theorien jenseits des Standardmodells ist.
Trotz dieser Schwächen im phänomenologischen Bereich ist die Supersymmetrie ein formal überaus ansprechendes Konzept, und es wäre interessant, das hier diskutierte Massenmodell in eine supersymmetrische Grand Unified-Theorie einzubetten und die daraus resultierenden Konsequenzen zu untersuchen. Auch eine Analyse des Modells hinsichtlich seiner Möglichkeiten zur Erklärung der $`CP`$-Verletzung und der Baryonasymmetrie des Universums wäre sicherlich sinnvoll.
Auf die generellen Probleme von nichtsupersymmetrischen Grand Unified-Theorien wurde in Kapitel 2 eingegangen. Das Hierarchieproblem und die Nichtberücksichtigung der Gravitation sind wohl die fundamentalsten Schwachpunkte dieser Modelle. Das Hierarchieproblem kann ebenso wie das Problem der Divergenzen zumindest prinzipiell durch eine bei Energien von etwa 1 TeV gebrochene globale Supersymmetrie gelöst werden. Bei solchen Modellen stellt sich aber die Frage, auf welche Weise die Supersymmetrie gebrochen wird. Sowohl die spontane Brechung als auch die explizite Brechung durch Lagrangedichteterme deuten letztendlich auf eine noch fundamentalere Theorie hin. Ein Kandidat für eine solche ist die Supergravitationstheorie, welche auf einer lokalen Supersymmetrie basiert und die Gravitationswechselwirkung mit einschließt. Auf Supergravitation basierende Modelle sind allerdings nicht renormierbar und können wiederum nur effektive Näherungen einer grundlegenden Theorie sein. Den zur Zeit zweifellos populärsten und vielleicht sogar einzigen ernstzunehmenden Ansatz dafür stellen die Superstringtheorien und das übergeordnete Modell der M-Theorie dar. Sie werden in zehn beziehungsweise elf Raumzeitdimensionen formuliert, was eine Kompaktifizierung der überzähligen Dimensionen notwendig macht, um vierdimensionale effektive Theorien und möglicherweise überprüfbare Voraussagen erhalten zu können. In der äußerst eingeschränkten Vorhersagekraft für den Bereich experimentell zugänglicher Energien besteht die Schwäche dieser Modelle.
Es ist bis heute keine wirklich fundamentale Theorie bekannt, welche die Eigenschaften der Fermionen erklärt. Aus diesem Grunde kommt der Untersuchung effektiver Massenmodelle nach wie vor große Bedeutung zu, da sie vielleicht einen Hinweis auf die zugrundeliegende Symmetrie der Natur liefern.
## Anhang A Gruppentheoretischer Anhang
Die in diesem Anhang zusammengefaßten Eigenschaften ausgewählter irreduzibler Darstellungen sind den Tabellen in entnommen; dort wird ferner eine detaillierte Beschreibung des für Anwendungen in der Elementarteilchenphysik relevanten gruppentheoretischen Formalismus gegeben.
### A.1 Casimir-Operator und Dynkin-Index
Die Generatoren einer Lie-Algebra erfüllen unabhängig von der expliziten Darstellung die Vertauschungsrelationen
$$[T_a,T_b]=if_{abc}T_c,$$
(A.1)
wobei die $`f_{abc}`$ die Strukturkonstanten der Algebra sind. Der quadratische Casimir-Operator $`\widehat{C}_2_aT_aT_a`$ kommutiert mit allen Generatoren und ist somit bezüglich jeder irreduziblen Darstellung $``$ der Liegruppe $`G`$ ein Vielfaches der Einheitsmatrix; der zugehörige Eigenwert wird mit $`C_2()`$ bezeichnet:
$$\widehat{C}_2()=C_2()$$
(A.2)
Der Dynkin-Index $`S_2()`$ einer irreduziblen Darstellung $``$ ist über
$$\text{Tr}_{}(T_aT_b)=S_2()\delta _{ab}$$
(A.3)
definiert. Zwischen $`C_2()`$ und $`S_2()`$ gilt die Beziehung
$$d(𝒢)S_2()=d()C_2(),$$
(A.4)
wobei $`d()`$ die Dimension von $``$ bezeichnet; $`𝒢`$ steht für die adjungierte Darstellung von $`G`$. Für die $`SU(N)`$-Gruppen gelten die Normierungen $`C_2(𝒢)=S_2(𝒢)=N`$ und $`S_2(N)=1/2`$, während für die $`U(1)_Y`$ $`C_2()=S_2()=Y^2`$ und $`C_2(𝒢)=S_2(𝒢)=0`$ ist.
In der folgenden Tabelle sind $`S_2()`$ und $`C_2()`$ einiger irreduzibler Darstellungen verschiedener Lie-Gruppen aufgelistet.
### A.2 Verzweigungsregeln für $`SO(10)`$-Darstellungen
$`SO(10)`$ $``$ $`SU(4)_CSU(2)_LSU(2)_R`$
$`\mathrm{𝟏𝟎}`$ $`=`$ $`(\mathbf{1,2,2})(\mathbf{6,1,1})`$ (A.5)
$`\mathrm{𝟏𝟔}`$ $`=`$ $`(\mathbf{4,2,1})(\overline{\mathrm{𝟒}}\mathbf{,1,2})`$ (A.6)
$`\mathrm{𝟒𝟓}`$ $`=`$ $`(\mathbf{15,1,1})(\mathbf{1,3,1})(\mathbf{1,1,3})(\mathbf{6,2,2})`$ (A.7)
$`\mathrm{𝟓𝟒}`$ $`=`$ $`(\mathbf{1,1,1})(\mathbf{1,3,3})(\mathrm{𝟐𝟎}^{}\mathbf{,1,1})(\mathbf{6,2,2})`$ (A.8)
$`\mathrm{𝟏𝟐𝟎}`$ $`=`$ $`(\mathbf{1,2,2})(\mathbf{10,1,1})(\overline{\mathrm{𝟏𝟎}}\mathbf{,1,1})(\mathbf{6,3,1})`$ (A.9)
$`(\mathbf{6,1,3})(\mathbf{15,2,2})`$
$`\mathrm{𝟏𝟐𝟔}`$ $`=`$ $`(\mathbf{6,1,1})(\overline{\mathrm{𝟏𝟎}}\mathbf{,3,1})(\mathbf{10,1,3})(\mathbf{15,2,2})`$ (A.10)
$`\mathrm{𝟐𝟏𝟎}`$ $`=`$ $`(\mathbf{1,1,1})(\mathbf{15,1,1})(\mathbf{6,2,2})(\mathbf{15,3,1})(\mathbf{15,1,3})`$ (A.11)
$`(\mathbf{10,2,2})(\overline{\mathrm{𝟏𝟎}}\mathbf{,2,2})`$
### A.3 Verzweigungsregeln für $`SU(4)`$-Darstellungen
$`SU(4)_C`$ $``$ $`SU(3)_CU(1)_{BL}`$
$`\mathrm{𝟒}`$ $`=`$ $`\mathrm{𝟏}_1\mathrm{𝟑}_{1/3}`$ (A.12)
$`\mathrm{𝟔}`$ $`=`$ $`\mathrm{𝟑}_{2/3}\overline{\mathrm{𝟑}}_{2/3}`$ (A.13)
$`\mathrm{𝟏𝟎}`$ $`=`$ $`\mathrm{𝟏}_2\mathrm{𝟑}_{2/3}\mathrm{𝟔}_{2/3}`$ (A.14)
$`\mathrm{𝟏𝟓}`$ $`=`$ $`\mathrm{𝟏}_0\mathrm{𝟑}_{4/3}\overline{\mathrm{𝟑}}_{4/3}\mathrm{𝟖}_0`$ (A.15)
$`\mathrm{𝟐𝟎}`$ $`=`$ $`\mathrm{𝟑}_{1/3}\overline{\mathrm{𝟑}}_{5/3}\overline{\mathrm{𝟔}}_{1/3}\mathrm{𝟖}_1`$ (A.16)
### A.4 Schwellenkorrektur-Koeffizienten
Hier ist zu berücksichtigen, daß die $`SO(10)`$-Darstellungen 210 und 54 als reell, die an der Massenerzeugung beteiligten Darstellungen 10, 120 und 126 dagegen als komplex angenommen werden. Hat die Darstellung einer Produktgruppe $`G_1G_2G_3`$ die Gestalt $`(_\mathrm{𝟏},_\mathrm{𝟐},_\mathrm{𝟑})`$, so sind die zugehörigen Schwellenkorrektur-Koeffizienten durch
$$(\lambda _1,\lambda _2,\lambda _3)=(d(_\mathrm{𝟐})d(_\mathrm{𝟑})S_2(_\mathrm{𝟏}),d(_\mathrm{𝟏})d(_\mathrm{𝟑})S_2(_\mathrm{𝟐}),d(_\mathrm{𝟏})d(_\mathrm{𝟐})S_2(_\mathrm{𝟑}))$$
(A.17)
gegeben. Bei komplexen Darstellungen kommt noch ein Faktor 2 hinzu.
Für die Berechnung der $`\lambda _{1Y}^I`$ ist die korrekte GUT-Normierung der Hyperladung verwendet worden, was die $`\sqrt{3/5}`$-Faktoren erklärt.
## Anhang B Renormierungsgruppengleichungen
In sind allgemeine Ausdrücke für die Renormierungsgruppengleichungen von Quantenfeldtheorien in Zweischleifenordnung entwickelt worden, aus denen sich auch die hier verwendeten Gleichungen herleiten lassen. Die Renormierungsgruppengleichungen für das Standardmodell sind in und die der See-Saw-Massenmatrix der Neutrinos in ausführlich analysiert worden.
### B.1 Standardmodell
Alle Gleichungen bis auf die der See-Saw-Massenmatrix sind in Zweischleifenordnung angegeben. Als Variable wird im folgenden $`t=\mathrm{ln}(\mu /M_Z)`$ verwendet, $`n_H`$ gibt die Anzahl der komplexen Higgs-Doubletts an (in dieser Arbeit wird stets $`n_H=1`$ angenommen). Die Definition der Yukawa-Kopplungsmatrizen lautet
$$𝐘_i(\mu )=\frac{\sqrt{2}}{\upsilon }𝐌_i(\mu )(i=u,d,e),$$
(B.1)
wobei $`𝐌_i(\mu )`$ die skalenabhängigen fermionischen Massenmatrizen bezeichnet; $`\upsilon /\sqrt{2}=174.1\text{GeV}`$ ist der Vakuumerwartungswert des Higgs-Feldes. Desweiteren werden die nachfolgend aufgeführten Definitionen benutzt:
$`𝐇_i`$ $`=`$ $`𝐘_i𝐘_i^{}(i=u,d,e)`$ (B.2)
$`Y_2`$ $`=`$ $`\text{Tr}\left(3𝐇_u+3𝐇_d+𝐇_e\right)`$ (B.3)
$`Y_4`$ $`=`$ $`\left({\displaystyle \frac{17}{20}}g_1^2+{\displaystyle \frac{9}{4}}g_2^2+8g_3^2\right)\text{Tr}\left(𝐇_u\right)`$ (B.4)
$`+`$ $`\left({\displaystyle \frac{1}{4}}g_1^2+{\displaystyle \frac{9}{4}}g_2^2+8g_3^2\right)\text{Tr}\left(𝐇_d\right)`$
$`+`$ $`\left({\displaystyle \frac{3}{4}}g_1^2+{\displaystyle \frac{3}{4}}g_2^2\right)\text{Tr}\left(𝐇_e\right)`$
$`H_2`$ $`=`$ $`\text{Tr}\left(3𝐇_u^2+3𝐇_d^2+𝐇_e^2\right)`$ (B.5)
$`\chi _4`$ $`=`$ $`{\displaystyle \frac{9}{4}}\text{Tr}\left(3𝐇_u^2+3𝐇_d^2+𝐇_e^2{\displaystyle \frac{2}{3}}𝐇_u𝐇_d\right)`$ (B.6)
#### B.1.1 Renormierungsgruppengleichungen für die Eichkopplungen
$`{\displaystyle \frac{d}{dt}}\alpha _1`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}\left(4{\displaystyle \frac{1}{10}}n_H\right)\alpha _1^2`$ (B.7)
$`{\displaystyle \frac{1}{8\pi ^2}}\left[\left({\displaystyle \frac{19}{5}}{\displaystyle \frac{9}{50}}n_H\right)\alpha _1+\left({\displaystyle \frac{9}{5}}{\displaystyle \frac{9}{10}}n_H\right)\alpha _2+\left({\displaystyle \frac{44}{5}}\right)\alpha _3\right]\alpha _1^2`$
$`{\displaystyle \frac{1}{32\pi ^3}}\text{Tr}\left({\displaystyle \frac{17}{10}}𝐇_u+{\displaystyle \frac{1}{2}}𝐇_d+{\displaystyle \frac{3}{2}}𝐇_e\right)\alpha _1^2`$
$`{\displaystyle \frac{d}{dt}}\alpha _2`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}\left({\displaystyle \frac{10}{3}}{\displaystyle \frac{1}{6}}n_H\right)\alpha _2^2`$ (B.8)
$`{\displaystyle \frac{1}{8\pi ^2}}\left[\left({\displaystyle \frac{3}{5}}{\displaystyle \frac{3}{10}}n_H\right)\alpha _1+\left({\displaystyle \frac{11}{3}}{\displaystyle \frac{13}{6}}n_H\right)\alpha _2+\left(12\right)\alpha _3\right]\alpha _2^2`$
$`{\displaystyle \frac{1}{32\pi ^3}}\text{Tr}\left({\displaystyle \frac{3}{2}}𝐇_u+{\displaystyle \frac{3}{2}}𝐇_d+{\displaystyle \frac{1}{2}}𝐇_e\right)\alpha _2^2`$
$`{\displaystyle \frac{d}{dt}}\alpha _3`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}\mathrm{\hspace{0.17em}7}\alpha _3^2`$ (B.9)
$`{\displaystyle \frac{1}{8\pi ^2}}\left[\left({\displaystyle \frac{11}{10}}\right)\alpha _1+\left({\displaystyle \frac{9}{2}}\right)\alpha _2+26\alpha _3\right]\alpha _3^2`$
$`{\displaystyle \frac{1}{32\pi ^3}}\text{Tr}\left(2𝐇_u+2𝐇_d\right)\alpha _3^2`$
#### B.1.2 Renormierungsgruppengleichungen für die Yukawa-Matrizen
$`{\displaystyle \frac{d}{dt}}𝐘_i`$ $`=`$ $`\left({\displaystyle \frac{1}{16\pi ^2}}\beta _i^{\left(1\right)}+{\displaystyle \frac{1}{(16\pi ^2)^2}}\beta _i^{\left(2\right)}\right)𝐘_i(i=u,d,e)`$ (B.10)
$`\beta _u^{\left(1\right)}`$ $`=`$ $`{\displaystyle \frac{3}{2}}\left(𝐇_u𝐇_d\right)+Y_2\left({\displaystyle \frac{17}{20}}g_1^2+{\displaystyle \frac{9}{4}}g_2^2+8g_3^2\right)`$ (B.11)
$`\beta _d^{\left(1\right)}`$ $`=`$ $`{\displaystyle \frac{3}{2}}\left(𝐇_d𝐇_u\right)+Y_2\left({\displaystyle \frac{1}{4}}g_1^2+{\displaystyle \frac{9}{4}}g_2^2+8g_3^2\right)`$ (B.12)
$`\beta _e^{\left(1\right)}`$ $`=`$ $`{\displaystyle \frac{3}{2}}𝐇_e+Y_2\left({\displaystyle \frac{9}{4}}g_1^2+{\displaystyle \frac{9}{4}}g_2^2\right)`$ (B.13)
$`\beta _u^{\left(2\right)}`$ $`=`$ $`{\displaystyle \frac{3}{2}}𝐇_u^2𝐇_u𝐇_d{\displaystyle \frac{1}{4}}𝐇_d𝐇_u+{\displaystyle \frac{11}{4}}𝐇_d^2+Y_2\left({\displaystyle \frac{5}{4}}𝐇_d{\displaystyle \frac{9}{4}}𝐇_u\right)\chi _4+{\displaystyle \frac{3}{2}}\lambda ^2`$ (B.14)
$`2\lambda \left(3𝐇_u+𝐇_d\right)+\left({\displaystyle \frac{223}{80}}g_1^2+{\displaystyle \frac{135}{16}}g_2^2+16g_3^2\right)𝐇_u\left({\displaystyle \frac{43}{80}}g_1^2{\displaystyle \frac{9}{16}}g_2^2+16g_3^2\right)𝐇_d`$
$`+{\displaystyle \frac{5}{2}}Y_4+{\displaystyle \frac{1187}{600}}g_1^4{\displaystyle \frac{9}{20}}g_1^2g_2^2+{\displaystyle \frac{19}{15}}g_1^2g_3^2{\displaystyle \frac{23}{4}}g_2^4+9g_2^2g_3^2108g_3^4`$
$`\beta _d^{\left(2\right)}`$ $`=`$ $`{\displaystyle \frac{3}{2}}𝐇_d^2𝐇_d𝐇_u{\displaystyle \frac{1}{4}}𝐇_u𝐇_d+{\displaystyle \frac{11}{4}}𝐇_u^2+Y_2\left({\displaystyle \frac{5}{4}}𝐇_u{\displaystyle \frac{9}{4}}𝐇_d\right)\chi _4+{\displaystyle \frac{3}{2}}\lambda ^2`$ (B.15)
$`2\lambda \left(3𝐇_d+𝐇_u\right)+\left({\displaystyle \frac{187}{80}}g_1^2+{\displaystyle \frac{135}{16}}g_2^2+16g_3^2\right)𝐇_d\left({\displaystyle \frac{79}{80}}g_1^2{\displaystyle \frac{9}{16}}g_2^2+16g_3^2\right)𝐇_u`$
$`+{\displaystyle \frac{5}{2}}Y_4{\displaystyle \frac{127}{600}}g_1^4{\displaystyle \frac{27}{20}}g_1^2g_2^2+{\displaystyle \frac{31}{15}}g_1^2g_3^2{\displaystyle \frac{23}{4}}g_2^4+9g_2^2g_3^2108g_3^4`$
$`\beta _e^{\left(2\right)}`$ $`=`$ $`{\displaystyle \frac{3}{2}}𝐇_e^2{\displaystyle \frac{9}{4}}Y_2𝐇_e\chi _4+{\displaystyle \frac{3}{2}}\lambda ^26\lambda 𝐇_e+\left({\displaystyle \frac{387}{80}}g_1^2+{\displaystyle \frac{135}{15}}g_2^2\right)𝐇_e+{\displaystyle \frac{5}{2}}Y_4`$ (B.16)
$`+{\displaystyle \frac{1371}{200}}g_1^4+{\displaystyle \frac{27}{20}}g_1^2g_2^2{\displaystyle \frac{23}{4}}g_2^4`$
#### B.1.3 Renormierungsgruppengleichung für die Higgs-Selbstkopplung
$`{\displaystyle \frac{d}{dt}}\lambda `$ $`=`$ $`\left({\displaystyle \frac{1}{16\pi ^2}}\beta _\lambda ^{\left(1\right)}+{\displaystyle \frac{1}{(16\pi ^2)^2}}\beta _\lambda ^{\left(2\right)}\right)\lambda `$ (B.17)
$`\beta _\lambda ^{\left(1\right)}`$ $`=`$ $`12\lambda ^2\left({\displaystyle \frac{9}{5}}g_1^2+9g_2^2\right)\lambda +{\displaystyle \frac{9}{4}}\left({\displaystyle \frac{3}{25}}g_1^4+{\displaystyle \frac{2}{5}}g_1^2g_2^2+g_2^4\right)+4Y_2\lambda 4H_2`$ (B.18)
$`\beta _\lambda ^{\left(2\right)}`$ $`=`$ $`78\lambda ^3+18\left({\displaystyle \frac{3}{5}}g_1^2+3g_2^2\right)\lambda ^2\left({\displaystyle \frac{2661}{100}}g_1^4{\displaystyle \frac{117}{20}}g_1^2g_2^2+{\displaystyle \frac{73}{8}}g_2^4\right)\lambda `$ (B.19)
$`{\displaystyle \frac{3411}{1000}}g_1^6{\displaystyle \frac{1677}{200}}g_1^4g_2^2{\displaystyle \frac{289}{40}}g_1^2g_2^4+{\displaystyle \frac{305}{8}}g_2^664g_3^2\text{Tr}\left(𝐇_u^2+𝐇_d^2\right)`$
$`{\displaystyle \frac{8}{5}}g_1^2\text{Tr}\left(2𝐇_u^2𝐇_d^2+3𝐇_e^2\right){\displaystyle \frac{3}{2}}g_2^4Y_2+10\lambda Y_4`$
$`+{\displaystyle \frac{3}{5}}g_1^2[({\displaystyle \frac{57}{10}}g_1^2+21g_2^2)\text{Tr}\left(𝐇_u\right)+({\displaystyle \frac{3}{2}}g_1^2+9g_2^2)\text{Tr}\left(𝐇_d\right)`$
$`+({\displaystyle \frac{15}{2}}g_1^2+11g_2^2)\text{Tr}\left(𝐇_e\right)]24\lambda ^2Y_2\lambda H_2+6\lambda \text{Tr}\left(𝐇_u𝐇_d\right))`$
$`+20\text{Tr}(3𝐇_u^3+3𝐇_d^3+𝐇_e^312\text{Tr}\left(𝐇_u^2𝐇_d\right)12\text{Tr}\left(𝐇_u𝐇_d^2\right)`$
#### B.1.4 Renormierungsgruppengleichung für die See-Saw-Massenmatrix
$`16\pi ^2{\displaystyle \frac{d}{dt}}𝐌_\nu `$ $`=`$ $`\left(2\lambda 3g_2^2+2Y_2\right)𝐌_\nu {\displaystyle \frac{1}{2}}\left(𝐌_\nu 𝐇_e+𝐇_e^T𝐌_\nu \right)`$ (B.20)
### B.2 $`SU(4)_CSU(2)_LSU(2)_R[D]`$-Modell
Die Renormierungsgruppengleichungen für die Eichkopplungen hängen vom Higgs-Teilcheninhalt des betrachteten Modells ab, das heißt von Anzahl und Art der Darstellungen, in welchen Higgs-Teilchen mit Massen der Größenordnung $`M_I`$ liegen. Folgende Variablen spezifizieren diesen Teilcheninhalt:
$`N_1`$ $`=`$ $`\text{Anzahl der}(\mathbf{1,2,2})_{10/120}`$
$`N_{15}`$ $`=`$ $`\text{Anzahl der}(\mathbf{15,2,2})_{120/126}`$
$`\mathrm{\Delta }_L`$ $`=`$ $`\text{Anzahl der}(\overline{\mathrm{𝟏𝟎}}\mathbf{,3,1})_{126}`$
$`\mathrm{\Delta }_R`$ $`=`$ $`\text{Anzahl der}(\mathbf{10,1,3})_{126}`$
Dann haben die Renormierungsgruppengleichungen für die Eichkopplungen in Zweischleifenordnung die folgende Gestalt, wobei Beiträge durch die Yukawakopplungen der Fermionen vernachlässigt werden:
$`{\displaystyle \frac{d}{dt}}\alpha _{4C}`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}\left({\displaystyle \frac{32}{3}}{\displaystyle \frac{16}{3}}N_{15}3\mathrm{\Delta }_L3\mathrm{\Delta }_R\right)\alpha _{4C}^2`$ (B.21)
$`{\displaystyle \frac{1}{8\pi ^2}}[({\displaystyle \frac{9}{2}}48N_{15}72\mathrm{\Delta }_L)\alpha _{2L}`$
$`+\left({\displaystyle \frac{9}{2}}48N_{15}72\mathrm{\Delta }_R\right)\alpha _{2R}`$
$`+({\displaystyle \frac{473}{6}}{\displaystyle \frac{896}{3}}N_{15}186\mathrm{\Delta }_L186\mathrm{\Delta }_R)\alpha _{4C}]\alpha ^2_{4C}`$
$`{\displaystyle \frac{d}{dt}}\alpha _{2L}`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}\left({\displaystyle \frac{10}{3}}{\displaystyle \frac{1}{3}}N_15N_{15}{\displaystyle \frac{20}{3}}\mathrm{\Delta }_L\right)\alpha _{2L}^2`$ (B.22)
$`{\displaystyle \frac{1}{8\pi ^2}}[({\displaystyle \frac{11}{3}}{\displaystyle \frac{13}{3}}N_165N_{15}{\displaystyle \frac{560}{3}}\mathrm{\Delta }_L)\alpha _{2L}`$
$`+\left(3N_145N_{15}\right)\alpha _{2R}`$
$`+({\displaystyle \frac{45}{2}}240N_{15}360\mathrm{\Delta }_L)\alpha _{4C}]\alpha ^2_{2L}`$
$`{\displaystyle \frac{d}{dt}}\alpha _{2R}`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}\left({\displaystyle \frac{10}{3}}{\displaystyle \frac{1}{3}}N_15N_{15}{\displaystyle \frac{20}{3}}\mathrm{\Delta }_R\right)\alpha _{2R}^2`$ (B.23)
$`{\displaystyle \frac{1}{8\pi ^2}}[(3N_145N_{15})\alpha _{2L}`$
$`+\left({\displaystyle \frac{11}{3}}{\displaystyle \frac{13}{3}}N_165N_{15}{\displaystyle \frac{560}{3}}\mathrm{\Delta }_R\right)\alpha _{2R}`$
$`+({\displaystyle \frac{45}{2}}240N_{15}360\mathrm{\Delta }_R)\alpha _{4C}]\alpha ^2_{2R}`$
## Anhang C Symmetriebrechungsskalen und Eichkopplungen
Die verwendeten Startwerte für die Integration der Renormierungsgruppengleichungen der Eichkopplungen sind in Tabelle C.1 zusammengefaßt.
### C.1 $`SU(4)_CSU(2)_LSU(2)_RD`$-Modell
Wie in Abschnitt 4.1.2 erläutert wurde, hängen der Wert von $`M_I`$ und somit die Eichkopplungen bei $`\mu =M_I`$ nicht vom Teilcheninhalt des Modells zwischen $`M_U`$ und $`M_I`$ ab; Tabelle C.2 gibt die Werte dieser Größen an.
Die vom Teilchenspektrum abhängigen Werte für $`M_U`$ und $`\alpha _U(M_U)`$ sind für die untersuchten Fälle in Tabelle C.3 aufgelistet.
### C.2 $`SU(4)_CSU(2)_LSU(2)_R`$-Modell
Im Gegensatz zum $`SU(4)_CSU(2)_LSU(2)_RD`$-Modell hängen in Theorien mit $`G_{\text{PS}}`$ als Symmetriegruppe auch $`M_I`$ und die Eichkopplungen bei der Skala $`\mu =M_I`$ vom Higgs-Spektrum ab. In Modellen mit $`N_{15}=3`$ und $`N_1=1,\mathrm{}\mathrm{,5}`$ sowie $`N_{15}=2`$ und $`N_1=\mathrm{1,2,3}`$ findet keine Vereinheitlichung statt, für die übrigen Fälle finden sich in Tabelle C.4 die Werte der relevanten Größen:
| Größe | Wert | Wert | Wert | Wert |
| --- | --- | --- | --- | --- |
| $`N_1`$ | 1 | 2 | 3 | 4 |
| $`N_{15}`$ | 1 | 1 | 1 | 1 |
| $`M_I`$ | $`1.7010^{10}`$ GeV | $`7.2510^{10}`$ GeV | $`1.9410^{11}`$ GeV | $`3.9910^{11}`$ GeV |
| $`\alpha _1(M_I)`$ | $`(46.85)^1`$ | $`(45.89)^1`$ | $`(45.24)^1`$ | $`(44.77)^1`$ |
| $`\alpha _2(M_I)`$ | $`(39.22)^1`$ | $`(39.94)^1`$ | $`(40.43)^1`$ | $`(40.79)^1`$ |
| $`\alpha _3(M_I)`$ | $`(29.95)^1`$ | $`(31.58)^1`$ | $`(32.69)^1`$ | $`(33.49)^1`$ |
| $`\alpha _{2R}(M_I)`$ | $`(58.22)^1`$ | $`(55.54)^1`$ | $`(53.72)^1`$ | $`(52.40)^1`$ |
| $`\alpha _{2L}(M_I)`$ | $`(39.22)^1`$ | $`(39.94)^1`$ | $`(40.43)^1`$ | $`(40.79)^1`$ |
| $`\alpha _{4C}(M_I)`$ | $`(29.98)^1`$ | $`(31.61)^1`$ | $`(32.71)^1`$ | $`(33.52)^1`$ |
| $`M_U`$ | $`6.0110^{16}`$ GeV | $`1.8310^{16}`$ GeV | $`8.0510^{15}`$ GeV | $`4.3710^{15}`$ GeV |
| $`\alpha _U(M_U)`$ | $`(32.37)^1`$ | $`(33.69)^1`$ | $`(34.57)^1`$ | $`(35.19)^1`$ |
## Anhang D Eigenschaften der untersuchten Lösungen
### D.1 Massenmatrizen, Massen und Mischungswinkel bei $`M_I`$
Tabelle D.1 gibt für die drei Beispielmodelle aus Abschnitt 4.4 die nichtverschwindenden Einträge der Massenmatrizen der geladenen Fermionen sowie deren Massen und Mischungswinkel bei $`M_I`$ an, wobei für die Mischungsmatrizen die Parametrisierung (1.34) benutzt wird. In Tabelle D.2 sind die nichtverschwindenden Elemente der Dirac-, Majorana- und See-Saw-Massenmatrizen sowie Massen und Mischungen der leichten Neutrinos zusammengefaßt.
### D.2 $`SO(10)`$-Higgs-Parameter bei $`M_I`$
Tabelle D.3 faßt für die betrachteten Lösungen die numerischen Werte der $`SO(10)`$-Higgs-Parameter aus (4.63-4.82) zusammen.
### D.3 Massenmatrizen, Massen und Mischungswinkel bei $`M_Z`$
Tabelle D.4 gibt für die drei untersuchten Modelle die Einträge der Massenmatrizen der geladenen Fermionen und der leichten Neutrinos sowie deren Massen und Mischungswinkel bei $`M_Z`$ an.
### D.4 CKM-Matrix und leptonische Mischungsmatrix bei $`M_Z`$
Im folgenden werden die CKM-Matrix V und die leptonische Mischungsmatrix U für die in Abschnitt 4.4 analysierten Modelle angegeben.
Modell 1:
$$𝐕=\left(\begin{array}{ccc}\hfill 0.9752& \hfill 0.2212& \hfill 0.0030\\ \hfill 0.2211& \hfill 0.9746& \hfill 0.0352\\ \hfill 0.0049& \hfill 0.0350& \hfill 0.9994\end{array}\right),𝐔=\left(\begin{array}{ccc}\hfill 0.891& \hfill 0.455& \hfill 0.005\\ \hfill 0.391& \hfill 0.771& \hfill 0.504\\ \hfill 0.233& \hfill 0.447& \hfill 0.864\end{array}\right)$$
(D.1)
Modell 2a:
$$𝐕=\left(\begin{array}{ccc}\hfill 0.9752& \hfill 0.2212& \hfill 0.0030\\ \hfill 0.2211& \hfill 0.9746& \hfill 0.0353\\ \hfill 0.0049& \hfill 0.0351& \hfill 0.9994\end{array}\right),𝐔=\left(\begin{array}{ccc}\hfill 0.999& \hfill 0.039& \hfill 0.012\\ \hfill 0.026& \hfill 0.844& \hfill 0.536\\ \hfill 0.031& \hfill 0.535& \hfill 0.844\end{array}\right)$$
(D.2)
Modell 2b:
$$𝐕=\left(\begin{array}{ccc}\hfill 0.9752& \hfill 0.2211& \hfill 0.0025\\ \hfill 0.2211& \hfill 0.9746& \hfill 0.0355\\ \hfill 0.0054& \hfill 0.0352& \hfill 0.9994\end{array}\right),𝐔=\left(\begin{array}{ccc}\hfill 0.999& \hfill 0.038& \hfill 0.012\\ \hfill 0.025& \hfill 0.833& \hfill 0.552\\ \hfill 0.031& \hfill 0.552& \hfill 0.834\end{array}\right)$$
(D.3)
## Anhang E Berechnung der Nukleonzerfallsraten
### E.1 $`A_i`$-Koeffizienten
In diesem Abschnitt wird die explizite Form der $`A_i`$-Koeffizienten dargestellt, welche in der effektiven Lagrangedichte (5.3) für die Nukleonenzerfälle auftreten. Hierbei sind $`\stackrel{~}{G}=g_U^2/2M_{X,Y}^2`$ und $`\stackrel{~}{G}^{}=g_U^2/2M_{X^{},Y^{}}^2`$, wobei $`M_{X,Y}^2=M_{X^{},Y^{}}^2M_U^2`$ angenommen wird.
$`A_1`$ $`=`$ $`\stackrel{~}{G}\left((𝐔_R)_{11}(𝐔_L)_{11}+(𝐔_R)_{21}(𝐔_L)_{21}+(𝐔_R)_{31}(𝐔_L)_{31}\right)`$
$`\left((𝐄_R)_{11}(𝐃_L)_{11}+(𝐄_R)_{21}(𝐃_L)_{21}+(𝐄_R)_{31}(𝐃_L)_{31}\right)`$
$`+`$ $`\stackrel{~}{G}\left((𝐔_R)_{11}(𝐃_L)_{11}+(𝐔_R)_{21}(𝐃_L)_{21}+(𝐔_R)_{31}(𝐃_L)_{31}\right)`$
$`\left((𝐄_R)_{11}(𝐔_L)_{11}+(𝐄_R)_{21}(𝐔_L)_{21}+(𝐄_R)_{31}(𝐔_L)_{31}\right)`$
$`A_2`$ $`=`$ $`\stackrel{~}{G}\left((𝐔_R)_{11}(𝐔_L)_{11}+(𝐔_R)_{21}(𝐔_L)_{21}+(𝐔_R)_{31}(𝐔_L)_{31}\right)`$
$`\left((𝐄_L)_{11}(𝐃_R)_{11}+(𝐄_L)_{21}(𝐃_R)_{21}+(𝐄_L)_{31}(𝐃_R)_{31}\right)`$
$`+`$ $`\stackrel{~}{G}^{}\left((𝐃_R)_{11}(𝐔_L)_{11}+(𝐃_R)_{21}(𝐔_L)_{21}+(𝐃_R)_{31}(𝐔_L)_{31}\right)`$
$`\left((𝐄_L)_{11}(𝐔_R)_{11}+(𝐄_L)_{21}(𝐔_R)_{21}+(𝐄_L)_{31}(𝐔_R)_{31}\right)`$
$`A_3`$ $`=`$ $`\stackrel{~}{G}\left((𝐔_R)_{11}(𝐔_L)_{11}+(𝐔_R)_{21}(𝐔_L)_{21}+(𝐔_R)_{31}(𝐔_L)_{31}\right)`$
$`\left((𝐄_R)_{12}(𝐃_L)_{11}+(𝐄_R)_{22}(𝐃_L)_{21}+(𝐄_R)_{32}(𝐃_L)_{31}\right)`$
$`+`$ $`\stackrel{~}{G}\left((𝐔_R)_{11}(𝐃_L)_{11}+(𝐔_R)_{21}(𝐃_L)_{21}+(𝐔_R)_{31}(𝐃_L)_{31}\right)`$
$`\left((𝐄_R)_{12}(𝐔_L)_{11}+(𝐄_R)_{22}(𝐔_L)_{21}+(𝐄_R)_{32}(𝐔_L)_{31}\right)`$
$`A_4`$ $`=`$ $`\stackrel{~}{G}\left((𝐔_R)_{11}(𝐔_L)_{11}+(𝐔_R)_{21}(𝐔_L)_{21}+(𝐔_R)_{31}(𝐔_L)_{31}\right)`$
$`\left((𝐄_L)_{12}(𝐃_R)_{11}+(𝐄_L)_{22}(𝐃_R)_{21}+(𝐄_L)_{32}(𝐃_R)_{31}\right)`$
$`+`$ $`\stackrel{~}{G}^{}\left((𝐃_R)_{11}(𝐔_L)_{11}+(𝐃_R)_{21}(𝐔_L)_{21}+(𝐃_R)_{31}(𝐔_L)_{31}\right)`$
$`\left((𝐄_L)_{12}(𝐔_R)_{11}+(𝐄_L)_{22}(𝐔_R)_{21}+(𝐄_L)_{32}(𝐔_R)_{31}\right)`$
$`A_5`$ $`=`$ $`\stackrel{~}{G}\left((𝐔_R)_{11}(𝐔_L)_{11}+(𝐔_R)_{21}(𝐔_L)_{21}+(𝐔_R)_{31}(𝐔_L)_{31}\right)`$
$`\left((𝐄_R)_{11}(𝐃_L)_{12}+(𝐄_R)_{21}(𝐃_L)_{22}+(𝐄_R)_{31}(𝐃_L)_{32}\right)`$
$`+`$ $`\stackrel{~}{G}\left((𝐔_R)_{11}(𝐃_L)_{12}+(𝐔_R)_{21}(𝐃_L)_{22}+(𝐔_R)_{31}(𝐃_L)_{32}\right)`$
$`\left((𝐄_R)_{11}(𝐔_L)_{11}+(𝐄_R)_{21}(𝐔_L)_{21}+(𝐄_R)_{31}(𝐔_L)_{31}\right)`$
$`A_6`$ $`=`$ $`\stackrel{~}{G}\left((𝐔_R)_{11}(𝐔_L)_{11}+(𝐔_R)_{21}(𝐔_L)_{21}+(𝐔_R)_{31}(𝐔_L)_{31}\right)`$
$`\left((𝐄_L)_{11}(𝐃_R)_{12}+(𝐄_L)_{21}(𝐃_R)_{22}+(𝐄_L)_{31}(𝐃_R)_{32}\right)`$
$`+`$ $`\stackrel{~}{G}^{}\left((𝐃_R)_{12}(𝐔_L)_{11}+(𝐃_R)_{22}(𝐔_L)_{21}+(𝐃_R)_{32}(𝐔_L)_{31}\right)`$
$`\left((𝐄_L)_{11}(𝐔_R)_{11}+(𝐄_L)_{21}(𝐔_R)_{21}+(𝐄_L)_{31}(𝐔_R)_{31}\right)`$
$`A_7`$ $`=`$ $`\stackrel{~}{G}\left((𝐔_R)_{11}(𝐔_L)_{11}+(𝐔_R)_{21}(𝐔_L)_{21}+(𝐔_R)_{31}(𝐔_L)_{31}\right)`$
$`\left((𝐄_R)_{12}(𝐃_L)_{12}+(𝐄_R)_{22}(𝐃_L)_{22}+(𝐄_R)_{32}(𝐃_L)_{32}\right)`$
$`+`$ $`\stackrel{~}{G}\left((𝐔_R)_{11}(𝐃_L)_{12}+(𝐔_R)_{21}(𝐃_L)_{22}+(𝐔_R)_{31}(𝐃_L)_{32}\right)`$
$`\left((𝐄_R)_{12}(𝐔_L)_{11}+(𝐄_R)_{22}(𝐔_L)_{21}+(𝐄_R)_{32}(𝐔_L)_{31}\right)`$
$`A_8`$ $`=`$ $`\stackrel{~}{G}\left((𝐔_R)_{11}(𝐔_L)_{11}+(𝐔_R)_{21}(𝐔_L)_{21}+(𝐔_R)_{31}(𝐔_L)_{31}\right)`$
$`\left((𝐄_L)_{12}(𝐃_R)_{12}+(𝐄_L)_{22}(𝐃_R)_{22}+(𝐄_L)_{32}(𝐃_R)_{32}\right)`$
$`+`$ $`\stackrel{~}{G}^{}\left((𝐃_R)_{12}(𝐔_L)_{11}+(𝐃_R)_{22}(𝐔_L)_{21}+(𝐃_R)_{32}(𝐔_L)_{31}\right)`$
$`\left((𝐄_L)_{12}(𝐔_R)_{11}+(𝐄_L)_{22}(𝐔_R)_{21}+(𝐄_L)_{32}(𝐔_R)_{31}\right)`$
$`A_9`$ $`=`$ $`\stackrel{~}{G}\left((𝐔_R)_{11}(𝐃_L)_{11}+(𝐔_R)_{21}(𝐃_L)_{21}+(𝐔_R)_{31}(𝐃_L)_{31}\right)`$
$`\left((𝐍_L)_{11}(𝐃_R)_{11}+(𝐍_L)_{21}(𝐃_R)_{21}+(𝐍_L)_{31}(𝐃_R)_{31}\right)`$
$`\stackrel{~}{G}^{}\left((𝐃_R)_{11}(𝐃_L)_{11}+(𝐃_R)_{21}(𝐃_L)_{21}+(𝐃_R)_{31}(𝐃_L)_{31}\right)`$
$`\left((𝐍_L)_{11}(𝐔_R)_{11}+(𝐍_L)_{21}(𝐔_R)_{21}+(𝐍_L)_{31}(𝐔_R)_{31}\right)`$
$`A_{10}`$ $`=`$ $`\stackrel{~}{G}\left((𝐔_R)_{11}(𝐃_L)_{11}+(𝐔_R)_{21}(𝐃_L)_{21}+(𝐔_R)_{31}(𝐃_L)_{31}\right)`$
$`\left((𝐍_L)_{12}(𝐃_R)_{11}+(𝐍_L)_{22}(𝐃_R)_{21}+(𝐍_L)_{32}(𝐃_R)_{31}\right)`$
$`\stackrel{~}{G}^{}\left((𝐃_R)_{11}(𝐃_L)_{11}+(𝐃_R)_{21}(𝐃_L)_{21}+(𝐃_R)_{31}(𝐃_L)_{31}\right)`$
$`\left((𝐍_L)_{12}(𝐔_R)_{11}+(𝐍_L)_{22}(𝐔_R)_{21}+(𝐍_L)_{32}(𝐔_R)_{31}\right)`$
$`A_{11}`$ $`=`$ $`\stackrel{~}{G}\left((𝐔_R)_{11}(𝐃_L)_{11}+(𝐔_R)_{21}(𝐃_L)_{21}+(𝐔_R)_{31}(𝐃_L)_{31}\right)`$
$`\left((𝐍_L)_{11}(𝐃_R)_{12}+(𝐍_L)_{21}(𝐃_R)_{22}+(𝐍_L)_{31}(𝐃_R)_{32}\right)`$
$`\stackrel{~}{G}^{}\left((𝐃_R)_{12}(𝐃_L)_{11}+(𝐃_R)_{22}(𝐃_L)_{21}+(𝐃_R)_{32}(𝐃_L)_{31}\right)`$
$`\left((𝐍_L)_{11}(𝐔_R)_{11}+(𝐍_L)_{21}(𝐔_R)_{21}+(𝐍_L)_{31}(𝐔_R)_{31}\right)`$
$`A_{12}`$ $`=`$ $`\stackrel{~}{G}\left((𝐔_R)_{11}(𝐃_L)_{11}+(𝐔_R)_{21}(𝐃_L)_{21}+(𝐔_R)_{31}(𝐃_L)_{31}\right)`$
$`\left((𝐍_L)_{12}(𝐃_R)_{12}+(𝐍_L)_{22}(𝐃_R)_{22}+(𝐍_L)_{32}(𝐃_R)_{32}\right)`$
$`\stackrel{~}{G}^{}\left((𝐃_R)_{12}(𝐃_L)_{11}+(𝐃_R)_{22}(𝐃_L)_{21}+(𝐃_R)_{32}(𝐃_L)_{31}\right)`$
$`\left((𝐍_L)_{12}(𝐔_R)_{11}+(𝐍_L)_{22}(𝐔_R)_{21}+(𝐍_L)_{32}(𝐔_R)_{31}\right)`$
$`A_{13}`$ $`=`$ $`\stackrel{~}{G}\left((𝐔_R)_{11}(𝐃_L)_{12}+(𝐔_R)_{21}(𝐃_L)_{22}+(𝐔_R)_{31}(𝐃_L)_{32}\right)`$
$`\left((𝐍_L)_{11}(𝐃_R)_{11}+(𝐍_L)_{21}(𝐃_R)_{21}+(𝐍_L)_{31}(𝐃_R)_{31}\right)`$
$`\stackrel{~}{G}^{}\left((𝐃_R)_{11}(𝐃_L)_{12}+(𝐃_R)_{21}(𝐃_L)_{22}+(𝐃_R)_{31}(𝐃_L)_{32}\right)`$
$`\left((𝐍_L)_{11}(𝐔_R)_{11}+(𝐍_L)_{21}(𝐔_R)_{21}+(𝐍_L)_{31}(𝐔_R)_{31}\right)`$
$`A_{14}`$ $`=`$ $`\stackrel{~}{G}\left((𝐔_R)_{11}(𝐃_L)_{12}+(𝐔_R)_{21}(𝐃_L)_{22}+(𝐔_R)_{31}(𝐃_L)_{32}\right)`$
$`\left((𝐍_L)_{12}(𝐃_R)_{11}+(𝐍_L)_{22}(𝐃_R)_{21}+(𝐍_L)_{32}(𝐃_R)_{31}\right)`$
$`\stackrel{~}{G}^{}\left((𝐃_R)_{11}(𝐃_L)_{12}+(𝐃_R)_{21}(𝐃_L)_{22}+(𝐃_R)_{31}(𝐃_L)_{32}\right)`$
$`\left((𝐍_L)_{12}(𝐔_R)_{11}+(𝐍_L)_{22}(𝐔_R)_{21}+(𝐍_L)_{32}(𝐔_R)_{31}\right)`$
$`A_{15}`$ $`=`$ $`\stackrel{~}{G}\left((𝐔_R)_{11}(𝐃_L)_{11}+(𝐔_R)_{21}(𝐃_L)_{21}+(𝐔_R)_{31}(𝐃_L)_{31}\right)`$
$`\left((𝐍_L)_{13}(𝐃_R)_{11}+(𝐍_L)_{23}(𝐃_R)_{21}+(𝐍_L)_{33}(𝐃_R)_{31}\right)`$
$`\stackrel{~}{G}^{}\left((𝐃_R)_{11}(𝐃_L)_{11}+(𝐃_R)_{21}(𝐃_L)_{21}+(𝐃_R)_{31}(𝐃_L)_{31}\right)`$
$`\left((𝐍_L)_{13}(𝐔_R)_{11}+(𝐍_L)_{23}(𝐔_R)_{21}+(𝐍_L)_{33}(𝐔_R)_{31}\right)`$
$`A_{16}`$ $`=`$ $`\stackrel{~}{G}\left((𝐔_R)_{11}(𝐃_L)_{11}+(𝐔_R)_{21}(𝐃_L)_{21}+(𝐔_R)_{31}(𝐃_L)_{31}\right)`$
$`\left((𝐍_L)_{13}(𝐃_R)_{12}+(𝐍_L)_{23}(𝐃_R)_{22}+(𝐍_L)_{33}(𝐃_R)_{32}\right)`$
$`\stackrel{~}{G}^{}\left((𝐃_R)_{12}(𝐃_L)_{11}+(𝐃_R)_{22}(𝐃_L)_{21}+(𝐃_R)_{32}(𝐃_L)_{31}\right)`$
$`\left((𝐍_L)_{13}(𝐔_R)_{11}+(𝐍_L)_{23}(𝐔_R)_{21}+(𝐍_L)_{33}(𝐔_R)_{31}\right)`$
$`A_{17}`$ $`=`$ $`\stackrel{~}{G}\left((𝐔_R)_{11}(𝐃_L)_{12}+(𝐔_R)_{21}(𝐃_L)_{22}+(𝐔_R)_{31}(𝐃_L)_{32}\right)`$
$`\left((𝐍_L)_{13}(𝐃_R)_{11}+(𝐍_L)_{23}(𝐃_R)_{21}+(𝐍_L)_{33}(𝐃_R)_{31}\right)`$
$`\stackrel{~}{G}^{}\left((𝐃_R)_{11}(𝐃_L)_{12}+(𝐃_R)_{21}(𝐃_L)_{22}+(𝐃_R)_{31}(𝐃_L)_{32}\right)`$
$`\left((𝐍_L)_{13}(𝐔_R)_{11}+(𝐍_L)_{23}(𝐔_R)_{21}+(𝐍_L)_{33}(𝐔_R)_{31}\right)`$
### E.2 Meson-Wellenfunktionen
Die Tabellen E.1 und E.2 geben die Spin-Flavour-Anteile der Wellenfunktionen für die pseudoskalaren beziehungsweise Vektormesonen an.
### E.3 Phasenraumfaktoren
Die Phasenraumfaktoren, welche die $`SU(6)`$-Spin-Flavour-Symmetrie des in Abschnitt 5.2 verwendeten nichtrelativistischen Quark-Modells brechen, sind durch
$$\rho _{p,n}=(1\chi _{p,n}^2)(1\chi _{p,n}^4)$$
(E.1)
gegeben, wobei $`\chi _{p,n}=m_{\text{Meson}}/m_{p,n}`$ der Quotient aus Mesonmasse und Masse des Protons beziehungsweise Neutrons ist ($`m_p=938.27231`$ MeV, $`m_n=939.56563`$ MeV ). Tabelle E.3 gibt die entsprechenden Werte für die relevanten Mesonen an.
### E.4 Übergangsamplituden
In diesem Abschnitt sind die Resultate für die hadronischen Übergangsamplituden und -wahrscheinlichkeiten der Nukleonenzerfälle zusammengefaßt. Tabelle E.4 zeigt die Amplituden der elementaren Zefallsprozesse des Protons und des Neutrons. Alle dort nicht aufgelisteten Amplituden lassen sich durch Symmetrieüberlegungen herleiten. So unterscheiden sich die Amplituden mit Antileptonen der zweiten und dritten Familie im Endzustand von denen mit $`e^+`$ und $`\nu _e^C`$ lediglich durch die Vorfaktoren $`A_i`$; diese lassen sich (5.3) entnehmen. Die Übergangsamplituden mit linkshändigen Antileptonen erhält man aus denjenigen mit rechtshändigen Antileptonen, indem man $`\overline{l}_R^C`$ durch $`\overline{l}_L^C`$ ersetzt und die Spineinstellungen $``$ und $``$ der Quarks vertauscht. Auch hier ändern sich lediglich die Vorfaktoren $`A_i`$ gemäß (5.3), die numerischen Werte der Amplituden in der dritten Spalte bleiben gleich. Tabelle E.5 listet die Übergangswahrscheinlichkeiten für Zerfallsprozesse mit physikalischen Teilchen in den Endzuständen auf.
### E.5 Experimentelle Grenzen für die Nukleonzerfallsraten
Tabelle E.6 zeigt die in angegebenen Untergrenzen für die inversen partiellen Zerfallsraten der Nukleonen aufgrund der experimentellen Nichtbeobachtung dieser Prozesse. Die aktuellsten Resultate des Super-Kamiokande-Experiments Nukleonenzerfälle betreffend sind in Tabelle 5.3 zu finden.
### E.6 Zerfallsraten der Nukleonen
In diesem Abschnitt sind in den Tabellen E.7-E.10 die numerischen Resultate für die partiellen und totalen Nukleonzerfallsraten zusammengefaßt. Neben den Ergebnissen für die drei untersuchten Modelle ist zum Vergleich auch der Fall verschwindender Fermionmischungen berücksichtigt worden. Während die Einträge 0.0 für partielle Raten $`<0.05\%`$ stehen, sind Zerfälle in die mit „—“ gekennzeichneten Kanäle in führender Ordnung nicht erlaubt. Es gilt $`\mathrm{\Gamma }=_j\mathrm{\Gamma }_j`$ sowie $`\tau =1/\mathrm{\Gamma }`$ und weiterhin in natürlichen Einheiten ($`\mathrm{}=c=1`$) 1 GeV$`=4.7710^{31}`$ yr<sup>-1</sup>.
## Schlußwort
An dieser Stelle möchte ich mich bei all denen bedanken, die mich während der letzten drei Jahre auf die eine oder andere Weise unterstützt haben.
Zunächst gilt mein Dank Prof. Dr. Yoav Achiman, der dieses Projekt angeregt und betreut hat. Er hat durch sein ständiges Interesse, die stete Bereitschaft zu ausgiebigen Diskussionen und zahlreiche physikalische Ideen wesentlich zum Gelingen dieser Arbeit beigetragen. Dank auch den übrigen Mitgliedern des Prüfungsausschusses, Prof. Dr. Peter Kroll als Zweitgutachter der Dissertation sowie Prof. Dr. Jürgen Drees und Dr. Peer Ueberholz.
Dr. Mina Parida war mir während seines Forschungsaufenthaltes in Wuppertal mit seinem Fachwissen eine sehr große Hilfe und stets bereit, sich mit meinem Fragen zu beschäftigen. Ebenso danke ich Prof. Dr. Franco Buccella, Dr. Ofelia Pisanti und Dr. Pietro Santorelli für die aufschlußreichen Diskussionen währen der Sommerschule auf Capri.
Dank auch allen Kollegen am Institut, insbesondere aber meinen Mitstreitern aus F 10.09. Marcus Richter, Dr. Armin Seyfried, Thorsten Struckmann und Dr. Jochen Viehoff verdanke ich eine überaus angenehme Arbeitsatmosphäre während der vergangenen drei Jahre. Auch die Hilfe von Norbert Eicker in EDV-Dingen sowie von Sabine Hoffmann und Anita Wied aus dem Sekretariat in allen organisatorischen Fragen ist nicht zu unterschätzen.
Meiner Familie und meinem Freundeskreis gebührt ebenfalls großer Dank. Meine Eltern haben mich während meines gesamten Studiums in jeder Weise tatkräftig unterstützt. Den Freunden, allen voran aber Christine Rompel, sei für ihre große Geduld und Aufmunterung vor allem in der Endphase meiner Promotion gedankt.
Der Deutschen Forschungsgemeinschaft schließlich bin ich für die finanzielle Unterstützung in Form eines Stipendiums im Rahmen des Graduiertenkollegs „Feldtheoretische und numerische Methoden der Elementarteilchen- und Statistischen Physik“ überaus dankbar.
Carsten Merten
Wuppertal, Dezember 1999 |
warning/0002/astro-ph0002430.html | ar5iv | text | # Neutrino Afterglows and Progenitors of Gamma-Ray Bursts
## 1 Introduction
The study of gamma-ray bursts (GRBs) has been revolutionized due to observations of multiwavelength afterglows in the past few years, but the nature of their progenitors remains unknown (for a review see Mészáros 1999). Two currently popular models for GRB progenitors are the mergers of compact objects (neutron stars or black holes) and the explosions of massive stars. In the former model, compact objects are expected to have significant spacial velocities so that their mergers would take place at many kiloparsecs outside their birthplaces. Thus, GRBs produced by this model would occur in the interstellar medium (ISM) with density $`n1\mathrm{cm}^3`$. Strong evidence for the massive star progenitor model has been recently discovered. GRB 980425 was probably associated with the relatively nearby Type Ic supernova (SN) 1998bw (Iwamoto et al. 1998; Kulkarni et al. 1998), and the supernova-like emission was also found in GRB 980326 (Bloom et al. 1999) and GRB 970228 (Reichart 1999; Galama et al. 2000). These observations show that some or possibly all long-duration GRBs arise from the core collapse of massive stars. It has been widely believed that GRBs associated with supernovae should unavoidably occur in the preburst stellar wind environment with mass density $`\rho R^2`$. If GRB emission is isotropic, the X-ray and optical afterglow in the wind case must decline more rapidly than in the ISM case, as studied analytically by Dai & Lu (1998), Mészáros, Rees & Wijers (1998), Panaitescu, Mészáros & Rees (1998) and Chevalier & Li (1999, 2000). Guided by this argument, Dai & Lu (1998) suggested for the first time that GRB 970616 is a possible wind interactor based on the rapid fading indicated by two X-ray flux measurements. Recently Chevalier & Li (1999) argued that GRB 980519 is an excellent wind interactor based on its X-ray, optical and radio data (although these observational data were shown analytically and numerically to be consistent with the dense medium model by Dai & Lu and Wang, Dai & Lu ). Furthermore, the afterglow data of some other bursts (e.g., GRB 970228, GRB 970508, GRB 980326 and GRB 980425) are also consistent with the wind environment model (Chevalier & Li 2000). These afterglows have been argued to be further evidence for massive stars as GRB progenitors.
In this paper we study neutrino afterglows from reverse shocks as a result of the interaction of relativistic fireballs with their surrounding wind matter by assuming that GRBs result from the explosions of massive stars. We find that the differential spectrum of neutrinos below $`3\times 10^{15}`$ eV is proportional to $`ϵ_\nu ^1`$ but the differential spectrum of neutrinos with energy from $`3\times 10^{15}`$ to $`3\times 10^{17}`$ eV steepens by one power of the energy. In addition, the expected flux of upward moving muons produced by neutrino interactions below a detector on the surface of the Earth is $`5`$ events per year per km<sup>2</sup> for typical parameters. We also find that this flux is $`10`$ times larger than estimated by Waxman & Bahcall (2000), who studied neutrino emission from reverse shocks produced by the interaction of fireballs with the interstellar medium (ISM).
The neutrinos are produced by $`\pi ^+`$ created in interactions between accelerated protons and synchrotron photons from accelerated electrons in a relativistic fireball. This neutrino emission during the GRB phase was studied in the internal shock models by Waxman & Bahcall (1997), Halzen (1998) and Rachen & Mészáros (1998). It was found that a fraction, $`0.1`$, of the fireball energy would be converted by photomeson production to a burst of neutrinos with typical energy of a few $`10^{14}`$ eV (but also see Vietri 1998). The property of such neutrino bursts is independent of whether the ambient matter is a stellar wind or a constant density medium. Observations of these bursts could be used to test the simultaneity of neutrino and photon arrival to an accuracy of $`1`$ s, the weak equivalence principle, and the vacuum neutrino oscillation theory (Waxman & Bahcall 1997).
The structure of this paper is as follows: In Section 2 we analyze reverse shocks produced during the interaction of ultra-relativistic fireballs with the surrounding wind matter and discuss the photon emission from these shocks. In Section 3 we investigate neutrino afterglow emission as a result of photo-meson interaction in the reverse shocks and in Section 4 we discuss the detectability of such afterglows. In the final section, several conclusions are drawn.
## 2 Shock Model and Photon Emission
We first assume that a relativistic GRB shell will interact with the surrounding stellar wind via two shocks: a reverse shock and a forward shock. The forward shock runs forward into the wind while the reverse shock sweeps up the shell material. The recently observed optical flash of GRB 990123 has been argued to come from a reverse shock (Akerlof et al. 1999; Sari & Piran 1999; Mészáros & Rees 1999). We believe that the reverse shock emission should be common for all GRBs. The shocked ambient and shell materials are in pressure balance and are separated by a contact discontinuity. We assume that these shocked materials are uniform and move together. Sari & Piran (1995) and Mitra (1998) considered the jump conditions for relativistic shocks and found the common Lorentz factor of the shocked materials measured in the unshocked medium frame,
$$\gamma =\frac{\overline{\xi }^{1/4}\mathrm{\Gamma }^{1/2}}{\sqrt{2}},$$
(1)
where $`\mathrm{\Gamma }`$ is the Lorentz factor of the unshocked shell measured in this frame and $`\overline{\xi }\rho _{\mathrm{sh}}/\rho _\mathrm{w}`$ is the ratio of proper mass densities of the unshocked shell and the unshocked ambient medium. The proper mass density of the ambient medium is expressed as
$$\rho _w=\frac{\dot{M}_w}{4\pi R^2V_w}AR^2,$$
(2)
where $`\dot{M}_w`$ and $`V_w`$ are the mass loss rate and wind velocity of the progenitor star, and $`A\dot{M}_w/(4\pi V_w)=10^5M_{}\mathrm{yr}^1/(4\pi \times 10^3\mathrm{km}\mathrm{s}^1)A_{}=5\times 10^{11}\mathrm{g}\mathrm{cm}^1A_{}`$ and $`R`$ is the radius of the shell in units of 1 cm (Chevalier & Li 1999, 2000). The proper mass density of the unshocked shell is given by
$$\rho _{\mathrm{sh}}=\frac{E_0}{4\pi R^2\mathrm{\Gamma }^2c^2\mathrm{\Delta }},$$
(3)
where $`E_0`$ and $`\mathrm{\Delta }`$ are the energy and the width (measured in the unshocked medium frame) of the initial shell. A typical value of $`A_{}1`$ for Wolf-Rayet stars is found from stellar mass-loss rates and wind velocities (Willis 1991; Chevalier & Li 2000 ). Since GRBs are believed to come from internal shocks, $`\mathrm{\Delta }`$ is approximately equal to the speed of light times the GRB durations and thus its typical value should be $`10`$ light seconds. The rapid variability of GRBs and their nonthermal spectra require that $`\mathrm{\Gamma }`$ be a few hundreds (Woods & Loeb 1995). From the observed fluences of some GRBs and their measured redshifts, $`E_0`$ is estimated to be between $`10^{52}`$ and $`10^{54}`$ ergs. A recent analysis by Freedman & Waxman (2000) also gives this estimate. Scaling the involved quantities with these typical values, we find
$$\overline{\xi }=655\frac{E_{53}}{A_{}\mathrm{\Delta }_{10}\mathrm{\Gamma }_{300}^2},$$
(4)
where $`E_{53}=E_0/10^{53}\mathrm{ergs}`$, $`\mathrm{\Gamma }_{300}=\mathrm{\Gamma }/300`$ and $`\mathrm{\Delta }_{10}`$ is in units of 10 light seconds. From equations (1) and (4), the Lorentz factor of the shocked shell material is rewritten as
$$\gamma =62\frac{E_{53}^{1/4}}{A_{}^{1/4}\mathrm{\Delta }_{10}^{1/4}}.$$
(5)
Following Sari & Piran (1995) and Mitra (1998), we further derive the Lorentz factor of the shocked shell material measured in the unshocked shell rest frame,
$$\gamma ^{}=\frac{\overline{\xi }^{1/4}\mathrm{\Gamma }^{1/2}}{\sqrt{2}}=\frac{1}{2}\frac{\mathrm{\Gamma }}{\gamma }=2.42\frac{A_{}^{1/4}\mathrm{\Delta }_{10}^{1/4}\mathrm{\Gamma }_{300}}{E_{53}^{1/4}},$$
(6)
which implies that the reverse shock is relativistic. After the reverse shock passes through the shell, the shock front disappears. Instead of maintaining a constant Lorentz factor (e.g., equation ), the shocked materials slow down with time based on the Blandford-McKee (1976) self-similar solution. In the following we discuss photon emission from the reverse shock.
Because of pressure balance across the contact discontinuity, the shocked shell material and the shocked wind material have not only the same bulk Lotentz factor but also the same internal energy density. According to relativistic shock jump conditions, we obtain the internal energy density of the shocked shell material,
$$e^{}=4\gamma ^2\rho _\mathrm{w}c^2=2\mathrm{\Gamma }\rho _\mathrm{w}c^2\overline{\xi }^{1/2}.$$
(7)
We assume that $`ϵ_e`$ and $`ϵ_B`$ are the fractions of the internal energy density (in the shocked material rest frame) that are carried by electrons and magnetic fields respectively. The minimum Lorentz factor of the electrons accelerated behind the reverse shock is approximated by $`\gamma _m(m_p/m_e)ϵ_e\gamma ^{}`$ (Waxman & Bahcall 2000), viz.,
$$\gamma _m445ϵ_{e,1}\left(\frac{A_{}^{1/4}\mathrm{\Delta }_{10}^{1/4}\mathrm{\Gamma }_{300}}{E_{53}^{1/4}}\right),$$
(8)
where $`ϵ_{e,1}=ϵ_e/0.1`$. Moreover, the magnetic field strength in the shocked shell material is given by
$$B^{}=(8\pi ϵ_Be^{})^{1/2}=\frac{(16\pi ϵ_B\mathrm{\Gamma }A\overline{\xi }^{1/2})^{1/2}}{t_b},$$
(9)
where $`t_b=R/c`$ is the time in the burster’s rest frame. Substituting the relation between this time and the observed time ($`t`$), $`t_b=2\gamma ^2t/(1+z)`$, into the above equation, and using equations (4) and (5), we further have
$$B^{}=3.8\times 10^3ϵ_{B,3}^{1/2}\left(\frac{1+z}{2}\right)\left(\frac{\mathrm{\Delta }_{10}^{1/4}A_{}^{3/4}}{E_{53}^{1/4}t_s}\right)\mathrm{G},$$
(10)
where $`ϵ_{B,3}=ϵ_B/10^3`$, $`z`$ is the redshift of the source and $`t_s=t/1\mathrm{s}`$. It should be noted that $`ϵ_e0.1`$ (Freedman & Waxman 2000), but $`ϵ_B`$ is highly uncertain and its reasonable value may be taken from $`10^2`$ to $`10^6`$. Several previous studies of GRB afterglows (e.g., Galama et al. 1999; Dai & Lu 1999, 2000; Wang, Dai & Lu 2000) gives $`ϵ_B10^410^6`$. A recent detailed study of the afterglows of GRBs 980703, 990123 and 990510 by Panaitescu & Kumar (2000) leads to $`ϵ_B10^310^4`$. In addition, Holland et al. (2000) find that $`ϵ_B`$ is as small as $`10^5`$. Therefore, we choose a typical value: $`ϵ_B10^3`$.
Let’s consider synchrotron radiation of the electrons accelerated behind the reverse shock. We first derive two characteristic frequencies of synchrotron photons: the typical frequency $`\nu _m`$ corresponding to the minimum electron Lorentz factor and the cooling frequency $`\nu _c`$. From equations (5), (8) and (10), we obtain the typical frequency in the observer frame,
$$\nu _m=\frac{\gamma \gamma _m^2}{1+z}\frac{eB^{}}{2\pi m_ec}=6.3\times 10^{16}ϵ_{e,1}^2ϵ_{B,3}^{1/2}\frac{A_{}\mathrm{\Delta }_{10}^{1/2}\mathrm{\Gamma }_{300}^2}{E_{53}^{1/2}t_s}\mathrm{Hz}$$
(11)
The cooling frequency corresponds to the cooling Lorentz factor $`\gamma _c`$, at which an electron is cooling on the dynamical expansion time. We believe that this Lorentz factor in the reverse shock will increase with time because of the cooling timescale $`B^2t^2`$. Initially, $`\gamma _c11`$. At this stage, a cooling electron may be non-relativistic and its kinetic energy $`E_e(1/2)\gamma m_ec^2\beta ^2`$. Thus, its cooling timescale due to cyclotron radiation, measured in the observer’s frame, can be estimated as $`t_0=(1+z)E_e/P_{\mathrm{cyc}}(\beta )`$, where $`P_{\mathrm{cyc}}(\beta )=(4/3)\sigma _Tc\gamma ^2\beta ^2B^2/(8\pi )`$ is the cyclotron radiation power (in the local observer frame) of an electron with velocity of $`\beta c`$ and with $`\sigma _T`$ being the Thomson scattering cross section. Using equation (10), we easily find
$$t_0=1.0ϵ_{B,3}\left(\frac{1+z}{2}\right)\frac{\mathrm{\Delta }_{10}^{1/4}A_{}^{5/4}}{E_{53}^{1/4}}\mathrm{s}.$$
(12)
Please note that $`t_0`$ is independent of $`\beta `$. This implies that at $`t<t_0`$ an electron accelerated to $`\gamma _m`$ in the magnetic field $`B^{}`$ will be able to cool to become non-relativistic, initially through synchrotron radiation and subseqently through cyclotron radiation, on the dynamical expansion time $`t`$. However, when $`t>t_0`$, the magnetic field ($`B^{}t^1`$) will become weaker and the cooling timescale due to cyclotron radiation must be longer than $`t`$, so an electron with $`\gamma _m`$ cannot cool to a non-relativistic velocity on time $`t`$. In this case, cyclotron radiation is no longer a cooling mechanism but it should be replaced by synchrotron radiation. Since for typical parameters $`t_01`$ s, which is much less than the durations of long GRBs from the collapse of massive stars, we will discuss the photon and neutrino emission from the reversely shocked matter at $`t>t_0`$ in the remaining text. According to Sari, Piran & Narayan (1998), the cooling Lorentz factor is defined by
$$\gamma \gamma _cm_ec^2=P_{\mathrm{syn}}(\gamma _c)t/(1+z),$$
(13)
where $`P_{\mathrm{syn}}(\gamma _c)=(4/3)\sigma _Tc\gamma ^2\gamma _c^2\beta ^2B^2/(8\pi )`$ is the synchrotron radiation power (in the local observer frame) of an electron with Lorentz factor of $`\gamma _c`$ (and $`\beta 1`$). Equation (13) leads to
$$\gamma _c=\frac{6\pi m_ec(1+z)}{\sigma _T\gamma B^2t}=2.0ϵ_{B,3}^1\left(\frac{1+z}{2}\right)^1\frac{E_{53}^{1/4}t_s}{\mathrm{\Delta }_{10}^{1/4}A_{}^{5/4}}.$$
(14)
It is clear that $`\gamma _c1`$ for $`t1`$ s, implying the cooling electrons are indeed relativistic. Using this equation, we further derive the cooling frequency in the observer frame,
$$\nu _c=\frac{\gamma \gamma _c^2}{1+z}\frac{eB^{}}{2\pi m_ec}=1.3\times 10^{12}ϵ_{B,3}^{3/2}\left(\frac{1+z}{2}\right)^2\frac{E_{53}^{1/2}t_s}{A_{}^2\mathrm{\Delta }_{10}^{1/2}}\mathrm{Hz}.$$
(15)
We can see from equations (11) and (15) that for typical parameters the cooling frequency is much lower than the typical frequency, indicating that all radiating electrons cool rapidly down to the cooling Lorentz factor. In other words, the shocked shell material is in the fast cooling regime. It is interesting to note that this conclusion has also been drawn by Chevalier & Li (2000). Therefore, the observed specific luminosity peaks at $`ϵ_ch\nu _c`$ rather than $`ϵ_mh\nu _m`$, with a peak value approximated by
$$L_{ϵ_c}=(2\pi \mathrm{})^1(1+z)\gamma N_eP_c^{}=2.7\times 10^{62}ϵ_{B,3}^{1/2}\left(\frac{1+z}{2}\right)\frac{E_{53}A_{}^{1/2}}{\mathrm{\Gamma }_{300}\mathrm{\Delta }_{10}}\mathrm{s}^1,$$
(16)
where $`N_e=E_0ct/[(1+z)\mathrm{\Gamma }m_pc^2\mathrm{\Delta }]`$ is the number of radiating electrons in the shocked shell region, and $`P_c^{}=m_ec^2\sigma _TB^{}/(3e)`$ is the power radiated per electron per unit frequency in the shocked shell rest frame.
We turn to derive the synchrotron self-absorption frequency of the reversely shocked matter. In the comoving frame of the shocked matter, the absorption coefficient for $`\nu ^{}>\nu _c^{}`$ is given by
$$\alpha _\nu ^{}^{}=\frac{\sqrt{3}e^3}{8\pi m_e}\left(\frac{3e}{2\pi m_e^3c^5}\right)^{p/2}(m_ec^2)^{p1}K\lambda B^{(p+2)/2}\nu ^{(p+4)/2}\mathrm{\Gamma }\left(\frac{3p+2}{12}\right)\mathrm{\Gamma }\left(\frac{3p+22}{12}\right),$$
(17)
where $`K=4(p1)\gamma ^{}\gamma _c(\rho _{\mathrm{sh}}/m_p)`$ and $`\lambda =(1/2)_0^\pi (\mathrm{sin}\alpha )^{(p+2)/2}\mathrm{sin}\alpha d\alpha `$ (Rybicki & Lightman 1979). In the present case, the electron distribution index $`p=2`$ because $`\nu _c\nu _m`$, and the width of the reverse shock $`R^{}=\gamma ct/(1+z)`$. Setting $`\tau (\nu _a^{})\alpha _{\nu _a^{}}^{}R^{}=1`$, we can derive the synchrotron self-absorption frequency in the obsever’s frame,
$$\nu _a=\frac{\gamma \nu _a^{}}{1+z}=2.2\times 10^{15}\left(\frac{1+z}{2}\right)^{1/3}\frac{A_{}^{1/6}E_{53}^{1/6}}{\mathrm{\Delta }_{10}^{1/6}\mathrm{\Gamma }_{300}^{1/3}t_s^{2/3}}\mathrm{Hz}.$$
(18)
Hence, $`\nu _c\nu _a\nu _m`$ for typical parameters.
We assume that the electrons behind the reverse shock follow a power law energy distribution, $`dn_e^{}/d\gamma _e\gamma _e^2`$ for $`\gamma _e\gamma _m`$ (Blandford & Eichler 1987). In this case, the synchrotron radiation spectrum is a broken power law:
$$L_{ϵ_\gamma }=\{\begin{array}{ccc}L_{ϵ_c}(ϵ_a/ϵ_c)^{1/2}(ϵ_\gamma /ϵ_a)^{5/2}\hfill & \mathrm{if}ϵ_\gamma ϵ_a\hfill & \\ L_{ϵ_c}(ϵ_\gamma /ϵ_c)^{1/2}\hfill & \mathrm{if}ϵ_aϵ_\gamma ϵ_m\hfill & \\ L_{ϵ_c}(ϵ_m/ϵ_c)^{1/2}(ϵ_\gamma /ϵ_m)^1\hfill & \mathrm{if}ϵ_\gamma ϵ_m,\hfill & \end{array}$$
(19)
where $`ϵ_a=h\nu _a`$. The protons behind the reverse shock are expected to be accelerated to the same power-law distribution as the electrons (with the maximum proton energy which will be estimated in the next section).
## 3 Neutrino Emission
For convenience, we hereafter denote the particle energy measured in the shocked shell rest frame with a prime, and the particle energy in the observer frame without prime, e.g., $`ϵ_\gamma =\gamma ϵ_\gamma ^{}/(1+z)`$. We now consider neutrino production in the wind case. Assuming $`n_\gamma ^{}(ϵ_\gamma ^{})dϵ_\gamma ^{}`$ to be the photon number density in the shocked shell rest frame and following Waxman & Bahcall (1997), we can write the fractional energy loss rate of a proton with energy $`ϵ_p^{}`$ due to pion production,
$$t_\pi ^{}_{}{}^{}1(ϵ_p^{})\frac{1}{ϵ_p^{}}\frac{dϵ_p^{}}{dt^{}}=\frac{1}{2\gamma _p^2}c_{ϵ_0}^{\mathrm{}}𝑑ϵ\sigma _\pi (ϵ)\xi (ϵ)ϵ_{ϵ/2\gamma _p}^{\mathrm{}}𝑑xx^2n^{}(x),$$
(20)
where $`\gamma _p=ϵ_p^{}/m_pc^2`$, $`\sigma _\pi (ϵ)`$ is the cross section for pion production for a photon with energy $`ϵ`$ in the proton rest frame, $`\xi (ϵ)`$ is the average fraction of energy lost to the pion, $`ϵ_0=0.15`$ GeV is the threshold energy, and the photon number density is related to the observed specific luminosity by $`n^{}(x)=L_{ϵ_\gamma }(\gamma x)/[4\pi R^2c(1+z)\gamma x]`$. Because of the $`\mathrm{\Delta }`$ resonance, we find that photo-meson production is dominated by the interaction with photons in the energy range $`ϵ_\gamma >ϵ_m`$. Considering the photon spectrum in this energy range, equation (20) leads to
$$t_\pi ^{}_{}{}^{}1(ϵ_p^{})=\frac{L_{ϵ_c}}{3\pi (1+z)R^2\gamma }\left(\frac{ϵ_c^{}}{ϵ_m^{}}\right)^{1/2}\left(\frac{\gamma _pϵ_m^{}}{ϵ_{\mathrm{peak}}}\right)\frac{\sigma _{\mathrm{peak}}\xi _{\mathrm{peak}}\mathrm{\Delta }ϵ}{ϵ_{\mathrm{peak}}},$$
(21)
where $`\sigma _{\mathrm{peak}}=5\times 10^{28}\mathrm{cm}^2`$ and $`\xi _{\mathrm{peak}}=0.2`$ at the resonance $`ϵ=ϵ_{\mathrm{peak}}=0.3`$ GeV, and $`\mathrm{\Delta }ϵ=0.2`$ GeV is the peak width. The fraction of energy loss of a proton with observed energy $`ϵ_p`$ by pion production, $`f_\pi (ϵ_p)`$, is defined by $`t_\pi ^{}_{}{}^{}1`$ times the expansion time of the shocked shell material ($`R/\gamma c`$). Thus, we have
$$f_\pi (ϵ_p)=2.0ϵ_{e,1}\left(\frac{1+z}{2}\right)^2\left(\frac{A_{}^{3/2}\mathrm{\Delta }_{10}^{1/2}}{E_{53}^{1/2}}\right)\left(\frac{ϵ_p}{10^{17}\mathrm{eV}}\right).$$
(22)
It is interesting to note that $`f_\pi (ϵ_p)`$ is independent of $`ϵ_B`$ and $`\mathrm{\Gamma }`$. Similarly to the cooling electron Lorentz factor defined by Sari et al. (1998) (see equation ), we can define the cooling proton energy $`ϵ_{p,c}`$ based on $`f_\pi (ϵ_{p,c})=1`$. According to equation (22), we find $`ϵ_{p,c}5\times 10^{16}`$ eV for typical parameters. This implies that the protons with energy $`ϵ_{p,c}`$ accelerated behind the reverse shock must lose almost all of their energy (viz., significant cooling) due to photo-meson interactions, but the protons with energy $`<ϵ_{p,c}`$ lose only a fraction ($`f_\pi `$) of their energy.
We now turn to discuss the neutrino spectrum. The photo-meson interactions include (1) production of $`\pi `$ mesons: $`p\gamma p+\pi ^0`$ and $`p\gamma n+\pi ^+`$, and (2) decay of $`\pi `$ mesons: $`\pi ^02\gamma `$ and $`\pi ^+\mu ^++\nu _\mu e^++\nu _e+\overline{\nu }_\mu +\nu _\mu `$. These processes produce neutrinos with energy $`0.05ϵ_p`$ (Waxman & Bahcall 1997). Since the protons with energy $`<ϵ_{p,c}`$ lose only a fraction ($`f_\pi ϵ_p`$) of their energy, the differential spectrum of neutrinos below the break energy $`3\times 10^{15}`$ eV is harder than the proton spectrum by one power of the energy. But since the protons with energy $`ϵ_{p,c}`$ accelerated behind the reverse shock must lose almost all of their energy, the neutrino spectrum above the break traces the proton spectrum. Therefore, if the differential spectrum of accelerated protons is assumed to be a power law form $`n(ϵ_p)ϵ_p^2`$, the differential neutrino spectrum is $`n(ϵ_\nu )ϵ_\nu ^1`$ below the break and $`n(ϵ_\nu )ϵ_\nu ^2`$ above the break.
The maximum energy of the resultant neutrinos is estimated as follows. This energy is determined by the maximum energy of the protons accelerated by the reverse shock. The typical Fermi acceleration time is $`t_a^{}=fR_L/c`$, where $`R_L=(1+z)ϵ_p/(\gamma eB^{})`$ is the Larmor radius and $`f`$ is of order unity (Hillas 1984). The requirement that this acceleration time is equal to the time for energy loss of protons ($`t_\pi ^{}`$) due to pion production leads to the maximum proton energy,
$$ϵ_{p,\mathrm{max}}=5.6\times 10^{18}f^{1/2}ϵ_{e,1}^{1/2}ϵ_{B,3}^{1/4}\left(\frac{1+z}{2}\right)^{1/2}\frac{E_{53}^{3/8}}{A_{}^{5/8}\mathrm{\Delta }_{10}^{3/8}}\mathrm{eV}.$$
(23)
From this equation, we can draw two conclusions: (1) For reasonable parameters, the maximum proton energy is $`6\times 10^{18}`$ eV, which is two orders of magnitude smaller than the maximum energy of the protons accelerated by the reverse shock in the ISM case (Waxman & Bahcall 2000). The physical conditions in the reverse shock for the ISM case imply that protons can be Fermi accelerated to $`10^{21}`$ eV (Waxman 1995; Vietri 1995; Milgrom & Usov 1995; see Waxman for recent reviews). (2) The maximum energy of neutrinos produced in the wind case is $`3\times 10^{17}`$ eV.
## 4 Detectability
We discuss the detectability of the neutrino afterglow emission in the wind case. Since the protons with energy $`5\times 10^{16}`$ eV must lose almost all of their energy due to photo-meson interactions, the present day neutrino energy density due to GRBs is approximately given by $`U_\nu =(1/2)(1/2)t_H\dot{E}`$, where the first factor $`1/2`$ accounts for the fact that about one half of the proton energy is lost to neutral pions which do not produce neutrinos, the second factor $`1/2`$ accounts for the fact that about one half of the energy in charged pions is transferred to $`\nu _\mu +\overline{\nu }_\mu `$, and $`t_H=10`$ Gyr is the Hubble time. Here we assume that $`\dot{E}=0.8\times 10^{44}\mathrm{erg}\mathrm{Mpc}^3\mathrm{yr}^1`$ is the production rate of GRB energy per unit volume (Waxman & Bahcall 2000). The neutrino flux is thus approximated by
$$\varphi _\nu \frac{c}{4\pi }\frac{U_\nu }{ϵ_\nu }4\times 10^{15}\left(\frac{ϵ_\nu }{3\times 10^{15}\mathrm{eV}}\right)^1\mathrm{cm}^2\mathrm{s}^1\mathrm{sr}^1.$$
(24)
The resulting high-energy neutrinos may be observed by detecting the Cherenkov light emitted by upward moving muons produced by neutrino interactions below a detector on the surface of the Earth (Gaisser, Halzen & Stanev 1995; Gandhi et al. 1998). Planned 1 km<sup>3</sup> detectors of high energy neutrinos include ICECUBE, ANTARES, NESTOR (Halzen 1999) and NuBE (Roy, Crawford & Trattner 1999). The probability that a neutrino could produce a high-energy muon in the detector is approximated by $`P_{\nu \mu }6\times 10^4(ϵ_\nu /3\times 10^{15}\mathrm{eV})^{0.5}`$. Using equation (24), we obtain the observed neutrino event rate in a detector,
$$N_{\mathrm{events}}=2\pi \varphi _\nu P_{\nu \mu }5\left(\frac{ϵ_\nu }{3\times 10^{15}\mathrm{eV}}\right)^{0.5}\mathrm{km}^2\mathrm{yr}^1.$$
(25)
This equation shows that a km<sup>2</sup> neutrino detector should detect each year about 5 neutrinos (with energy of $`3\times 10^{15}`$ eV) correlated with GRBs. For a GRB, its neutrino emission from the reverse shock in the wind case should be delayed to a few seconds after the main burst. Waxman & Bahcall (2000) have found $`f_\pi 0.1`$ for neutrino emission from reverse shocks in the ISM case (where the typical energy of neutrinos is $`3\times 10^{17}`$ eV). Using the same expression of $`P_{\nu \mu }`$, we have re-derived their neutrino event rate in a detector and obtained $`N_{\mathrm{events}}0.5\mathrm{km}^2\mathrm{yr}^1`$, which is smaller than our event rate by a factor of $`10`$.
## 5 Discussion and Conclusions
Neutrino bursts during the GRB phase were studied in the internal shock models by Waxman & Bahcall (1997) and Halzen (1998) who found that the neutrino event rate in a detector (mainly neutrinos with typical energy of a few $`10^{14}`$ eV) is $`26`$ events per year per km<sup>2</sup>, which is larger than our event rate by a factor of $`5`$. Compared with the analytical result of Waxman & Bahcall (2000), our discussions on neutrino afterglows in the wind case can lead to the following conclusions: (1) The protons with energy $`5\times 10^{16}`$ eV must lose almost all of their energy due to photo-meson interactions and thus the neutrino afterglow emission in the wind case is dominated by neutrinos with energy $`3\times 10^{15}`$ eV. (2) The maximum energy of the protons accelerated behind the reverse shock in the wind case is $`6\times 10^{18}`$ eV, so ultrahigh energy cosmic rays cannot be produced in this case. In addition, the maximum neutrino energy is $`3\times 10^{17}`$ eV. (3) The neutrino differential spectrum below $`3\times 10^{15}`$ eV is proportional to $`ϵ_\nu ^1`$ but the spectrum between $`3\times 10^{15}`$ and $`3\times 10^{17}`$ eV steepens by one power of the energy. (4) The observed neutrino event rate in the wind case is $`10`$ times larger than the one in the ISM case.
If GRB emission is isotropic, the optical afterglow in the wind case must decline more steeply than in the ISM case. This has been suggested as a plausible way of distinguishing between the GRB progenitor models (Chevalier & Li 1999, 2000). It is widely believed that GRBs may come from jets (Kulkarni et al. 1999; Castro-Tirado et al. 1999). As argued by Rhoads (1999) and Sari, Piran & Halpern (1999), the optical afterglow from a jet is likely to decay more rapidly at late times than at the early stage due to the lateral spreading effect. If this effect is true, however, both ISM and wind cases should show the same emission feature during the lateral spreading phase, and in particular on a timescale of days the wind density is similar to typical ISM densities so that an interaction with the wind would give results that are not too different from the ISM case (Chevalier & Li 2000; Livio & Waxman 1999). If GRBs are beamed, thus, their optical afterglow emission could not be used to discriminate the massive star progenitor model from the compact binary progenitor model. However, the neutrino afterglow emission discussed here is independent of whether the fireballs are isotropic or highly collimated. Therefore, neutrino afterglows, if detected, may be used to distinguish between GRB progenitor models based on differential spectra of observed neutrinos and their event rates in a detector.
What we want to point out is that the above conclusions are drawn by considering typical values of the parameters entering the fireball shock model. In fact, these parameters may have respective distributions. It is interesting to note that such distributions may lead to an event rate larger than estimated in this paper (Halzen & Hooper 1999).
We would like to thank Prof. F. Halzen and the referee for their valuable comments that enabled us to improve our paper, Y. F. Huang for reading carefully this paper, and X. Y. Wang for useful discussions. This work was supported partially by the National Natural Science Foundation of China (grants 19825109 and 19773007), and by the National Project on Fundamental Researches. |
warning/0002/math0002038.html | ar5iv | text | # Three Bimodules for Mansfield’s Imprimitivity Theorem
## 1. Introduction
C. K. Ng has recently observed that an abstract Morita equivalence between a restricted coaction crossed product $`A\times _{\delta |}G/N`$ and the iterated dual action crossed product $`A\times _\delta G\times _{\widehat{\delta }}N`$ can be pieced together from Green’s imprimitivity theorem \[5, Theorem 6\] and Katayama duality \[8, Theorem 8\], thus giving a relatively non-technical, nonconstructive proof of Mansfield’s imprimitivity theorem \[10, Theorem 27\]. However, in applications — especially those concerning induced representations — it is often necessary to work with an explicit bimodule. Because Morita equivalence relations are composed with one another by tensoring the corresponding imprimitivity bimodules together, Ng’s transitivity argument does implicitly provide a bimodule. Thus the natural question arises as to whether Ng’s bimodule is in fact isomorphic to Mansfield’s.
In more detail: Ng considers a nondegenerate reduced coaction $`\delta `$ of a locally compact group $`G`$ on a $`C^{}`$-algebra $`A`$, and a closed normal amenable subgroup $`N`$ of $`G`$. An application of Green’s theorem to the dual action $`(A\times _\delta G,G,\widehat{\delta })`$ gives an $`A\times _\delta G\times _{\widehat{\delta }}G\times _{\widehat{\widehat{\delta }}|}G/NA\times _\delta G\times _{\widehat{\delta }|}N`$ imprimitivity bimodule $`X_N^G(A\times _\delta G)`$. Moreover, looking closely at Katayama’s duality theorem, one can derive an isomorphism $`\mathrm{\Theta }`$ of $`A\times _\delta G\times _{\widehat{\delta }}G\times _{\widehat{\widehat{\delta }}|}G/N`$ onto $`(A\times _{\delta |}G/N)𝒦(L^2(G))`$, and this latter algebra is Morita equivalent to $`A\times _{\delta |}G/N`$ via the bimodule $`(A\times _{\delta |}G/N)L^2(G)`$. Implicitly, then, Ng’s $`A\times _\delta G\times _{\widehat{\delta }|}NA\times _{\delta |}G/N`$ imprimitivity bimodule is the tensor product
$$\stackrel{~}{X_N^G}(A\times _\delta G)_\mathrm{\Theta }\left((A\times _{\delta |}G/N)L^2(G)\right).$$
(Here the tilde denotes the reverse bimodule.) Let $`Y_{G/N}^G(A)`$ denote Mansfield’s imprimitivity bimodule. Then the question in question is precisely whether
(1.1)
$$X_N^G(A\times _\delta G)_{A\times G\times N}Y_{G/N}^G(A)(A\times _{\delta |}G/N)L^2(G)$$
as imprimitivity bimodules. In other words, modulo crossed-product duality, are Green and Mansfield induction inverses of one another?
This and related questions concerning actions, twisted actions, and twisted coactions are addressed in the present paper in the context of discrete groups and full coactions. Our approach is to exploit the natural equivariance of the Mansfield and Green bimodules. For instance, in Section 3, we consider a maximal discrete coaction $`(A,G,\delta )`$ and Mansfield’s $`A\times _\delta G\times _{\widehat{\delta }}GA`$ imprimitivity bimodule $`Y_{G/G}^G(A)`$, which carries a $`\widehat{\widehat{\delta }}`$$`\delta `$ compatible bimodule coaction $`\delta ^Y`$. If in addition $`N`$ is a normal subgroup of $`G`$, Theorem 3.1 states that
(1.2)
$$X_N^G(A\times _\delta G)_{A\times G\times N}Y_{G/N}^G(A)Y_{G/G}^G(A)\times _{\delta ^Y|}G/N$$
as $`A\times _\delta G\times _{\widehat{\delta }}G\times _{\widehat{\widehat{\delta }}|}G/NA\times _{\delta |}G/N`$ imprimitivity bimodules. The only difference between (1.1) and (1.2) is the bimodule on the right-hand sides; but these turn out to be isomorphic (see the proof of Theorem 7.1). Thus Mansfield induction of representations from $`\mathrm{Rep}A\times _{\delta |}G/N`$ to $`\mathrm{Rep}A\times _\delta G\times _{\widehat{\delta }}N`$ via $`Y_{G/N}^G(A)`$ can be “undone” by Green induction via $`X_N^G(A\times _\delta G)`$ followed by Katayama duality to get from $`\mathrm{Rep}A\times _\delta G\times _{\widehat{\delta }}G\times _{\widehat{\widehat{\delta }}|}G/N`$ back to $`\mathrm{Rep}A\times _{\delta |}G/N`$. In this sense we can usefully view Mansfield and Green induction as inverse to one another. In Section 4, we obtain results dual to those of Section 3, starting with an action instead of a coaction. In Section 5 we show that the results of Section 3 pass to twisted coactions, and in Section 6 we round out this square of ideas with a set of results for twisted actions.
In Section 7 we return to the comparison between Ng’s bimodule and Mansfield’s. This is done by first establishing (1.1) for full coactions and discrete groups, and then dropping it down to reduced coactions. In general, we feel that this approach — establishing results first for full coactions, and later passing to quotients if results for reduced coactions are desired — is more efficient and cleaner conceptually than working with reduced coactions directly. On the other hand, we work with discrete groups simply to avoid many of the technicalities associated with coactions of general locally compact groups; also, this is the only context in which we have induced algebras for coactions, which appear in Section 6. There is no reason to believe that the other results in this paper will not hold in the general case. In fact, Theorems 3.1 and 4.1 will appear in for locally compact groups, but only as a product of the extensive machinery developed therein. The proofs here are considerably more direct, and more instructive.
## 2. Preliminaries
In this preliminary section we collect the formulas relevant to crossed products and imprimitivity theorems involving actions and coactions of *discrete* groups. Because the groups are discrete, the theory acquires quite an algebraic flavor, and to take full advantage of this we translate the standard machinery involving locally compact groups to our context.
### Coactions
Let $`\delta :AAC^{}(G)`$ be a coaction of a discrete group $`G`$ on a $`C^{}`$-algebra $`A`$. Because $`G`$ is discrete, the spectral subspaces $`\{A_s:sG\}`$ of $`\delta `$ densely span $`A`$, and the union $`𝒜:=_{sG}A_s`$ <sup>1</sup><sup>1</sup>1more precisely, the disjoint union, but this abuse will cause no harm, since the spectral subspaces are linearly independent forms a Fell bundle over $`G`$. As shown in , the coaction on $`A`$ sits “between” a “maximal” coaction on the full cross-sectional algebra $`C^{}(𝒜)`$ and a “minimal” coaction on the reduced cross-sectional algebra $`C_r^{}(𝒜)`$. Although the crossed product $`A\times _\delta G`$ and the covariant representations of the coaction cannot distinguish among the various possibilities between the extremes $`C^{}(𝒜)`$ and $`C_r^{}(𝒜)`$, some other constructions can. In particular, the imprimitivity theorems we need require the full cross-sectional algebra. Therefore, we *assume throughout that $`A=C^{}(𝒜)`$*, and we call the coaction $`\delta `$ *maximal* in this case.
The crossed product $`A\times _\delta G`$ is densely spanned by the Cartesian product $`𝒜\times G`$, where $`A_s\times \{t\}`$ has the obvious vector space structure for all $`s,tG`$, and the multiplication and involution are given on the generators by
$`(a_r,s)(b_t,u)`$ $`=(a_rb_t,u)\text{if }s=tu\text{ (and }0\text{ if not)}`$
$`(a_s,t)^{}`$ $`=(a_s^{},st).`$
Since the coaction $`\delta `$ is maximal, $`A\times _\delta G`$ is the enveloping $`C^{}`$-algebra of the linear span of the generators; that is, any operation-preserving mapping of the generators into a $`C^{}`$-algebra $`C`$ extends uniquely to a homomorphism $`A\times _\delta GC`$. If $`(B,G,ϵ)`$ is another coaction and $`\varphi :AB`$ is an equivariant homomorphism (equivalently, $`\varphi (A_s)B_s`$ for each $`sG`$), then we can “integrate up” to get a homomorphism $`\varphi \times G:A\times _\delta GB\times _ϵG`$ defined on the generators by
$$(\varphi \times G)(a_s,t)=(\varphi (a_s),t).$$
The dual action $`\widehat{\delta }`$ of $`G`$ on $`A\times _\delta G`$ is given on the generators by
$$\widehat{\delta }_t(a_r,s)=(a_r,st^1).$$
If $`N`$ is a normal subgroup of $`G`$, the coaction $`\delta `$ restricts to a maximal coaction $`\delta |`$ of $`G/N`$ on $`A`$, and the crossed product $`A\times _{\delta |}G/N`$ is densely spanned by $`𝒜\times G/N`$, with operations
$`(a_r,sN)(b_t,uN)`$ $`=(a_rb_t,uN)\text{if }sN=tuN\text{ (and }0\text{ if not)}`$
$`(a_s,tN)^{}`$ $`=(a_s^{},stN),`$
and maps nondegenerately into $`M(A\times _\delta G)`$ by
$$(a_s,tN)\underset{nN}{}(a_s,tn)\text{(strictly convergent). }$$
There is a *decomposition* coaction $`\delta ^{\mathrm{dec}}`$ of $`G`$ on the restricted crossed product $`A\times _{\delta |}G/N`$, given on the generators by
$$\delta ^{\mathrm{dec}}(a_s,tN)=(a_s,tN)s.$$
A coaction $`(A,G,\delta )`$ is *twisted* over $`G/N`$ (see ) if there is an orthogonal family $`\{p_{tN}:tNG/N\}`$ of projections in $`M(A_e)`$ which sum strictly to $`1`$ in $`M(A_e)`$, and such that
$$a_sp_{tN}=p_{stN}a_s\text{for all }s,tG\text{}$$
The twisted crossed product $`A\times _{\delta ,G/N}G`$ is the quotient of $`A\times _\delta G`$ by the ideal generated by differences of the form
$$(a_s,t)(a_sp_{tN},t).$$
We denote the quotient map by $`q:A\times _\delta GA\times _{\delta ,G/N}G`$, and we write
$$[a_s,t]:=q(a_s,t).$$
The ideal $`\mathrm{ker}q`$ is invariant under the restriction $`\widehat{\delta }|_N`$, and we denote the corresponding action of $`N`$ on $`A\times _{\delta ,G/N}G`$ by $`\stackrel{~}{\delta }`$. We also define the “restriction” $`q|`$ by the commutative diagram
and we write
$$[a_s,tN]:=q|(a_s,tN).$$
It is shown in that
$$[a_s,tN]a_sp_{tN}$$
extends to an isomorphism $`q(A\times _{\delta |}G/N)A`$.
Let $`ϵ`$ be a maximal coaction of the quotient $`G/N`$ on $`A`$. It is shown in that there is an *induced* maximal coaction $`(\mathrm{Ind}A,G,\mathrm{Ind}ϵ)`$ with spectral subspaces
$$(\mathrm{Ind}A)_s=A_{sN}\times \{s\},$$
and where the generators have the coordinate-wise operations
$$(a_{sN},s)(b_{tN},t)=(a_{sN}b_{tN},st)\text{and}(a_{sN},s)^{}=(a_{sN}^{},s^1).$$
### Actions
Let $`\alpha :G\mathrm{Aut}B`$ be an action of the discrete group $`G`$ on a $`C^{}`$-algebra $`B`$. The crossed product $`B\times _\alpha G`$ is densely spanned by the Cartesian product $`B\times G`$, where $`B\times \{s\}`$ has the obvious vector space structure for all $`sG`$, and the multiplication and involution are given on the generators by
$$(a,r)(b,s)=(a\alpha _r(b),rs)\text{and}(a,r)^{}=(\alpha _{r^1}(a^{}),r^1).$$
Again, $`B\times _\alpha G`$ is the enveloping $`C^{}`$-algebra of the span of the generators, and if $`(C,G,\beta )`$ is another action and $`\varphi :BC`$ is an equivariant homomorphism, then we can integrate up to get a homomorphism $`\varphi \times G:B\times _\alpha GC\times _\beta G`$ defined on the generators by
$$(\varphi \times G)(b,s)=(\varphi (b),s).$$
The dual coaction $`\widehat{\alpha }`$ of $`G`$ on $`B\times _\alpha G`$ is given on the generators by
$$\widehat{\alpha }(a,r)=(a,r)r.$$
The decomposition action $`\alpha ^{\mathrm{dec}}`$ of $`G`$ on the restricted crossed product $`B\times _{\alpha |}N`$ is given on the generators by
$$\alpha _r^{\mathrm{dec}}(a,n)=(\alpha _r(a),rnr^1).$$
An action $`(B,G,\alpha )`$ is twisted over $`N`$ (see ) if there is a unitary homomorphism $`nu_n:NM(B)`$ such that
$$\alpha _s(u_n)=u_{sns^1}\text{and}\alpha _n(b)=u_nbu_n^{}.$$
The twisted crossed product $`B\times _{\alpha ,N}G`$ is the quotient of $`B\times _\alpha G`$ by the ideal generated by differences of the form
$$(bu_n,s)(b,ns).$$
We denote the quotient map by $`q:B\times _\alpha GB\times _{\alpha ,N}G`$, and we write
$$[b,s]:=q(b,s).$$
We denote the dual coaction of $`G/N`$ on $`B\times _{\alpha ,N}G`$ by $`\stackrel{~}{\alpha }`$. We also define the “restriction” $`q|`$ by the commutative diagram
and we write
$$[b,n]:=q|(b,n).$$
It is shown in that
$$[b,n]bu_n$$
extends to an isomorphism $`q(B\times _{\alpha |}N)B`$.
Let $`\beta `$ be an action of the normal subgroup $`N`$ on $`B`$. We identify the induced algebra $`\mathrm{Ind}B`$ as the $`c_0`$-section algebra of a $`C^{}`$-bundle $`G\times _NB`$ over $`G/N`$. Specifically, $`N`$ acts diagonally on the trivial $`C^{}`$-bundle $`G\times BG`$ by $`n(s,b):=(sn^1,nb)`$, and the associated orbit space $`G\times _NB`$ has a natural $`C^{}`$-bundle structure over $`G/N`$: we denote the orbit of a pair $`(s,b)`$ by $`[s,b]`$, and the fiber over $`sN`$ is $`\{[sn,b]:nN,bB\}`$. The induced action $`(\mathrm{Ind}B,G,\mathrm{Ind}\beta )`$ is given on the generators by $`\mathrm{Ind}\beta _s([t,b])=[st,b]`$.
### Commutative diagrams of bimodules
Most of this paper concerns imprimitivity bimodules, but a few times (in the Sections 5 and 6) we will need the following more general concept: a *right-Hilbert* $`AB`$ bimodule is a Hilbert $`B`$-module $`X`$ equipped with a left $`A`$-module action by adjointable maps (which are automatically bounded and $`B`$-linear). We will *assume throughout* that the right inner product is *full* and the left action is *nondegenerate*. For example, a surjective homomorphism $`\varphi :AB`$ determines a right-Hilbert bimodule $`{}_{A}{}^{}B_{B}^{}`$ with the obvious right $`B`$-module action, left $`A`$-module action $`ab:=\varphi (a)b`$, and right inner product $`b,c_B:=b^{}c`$. Moreover, in this situation any right-Hilbert $`BC`$ bimodule $`X`$ can also be regarded as a right-Hilbert $`AC`$ bimodule with left $`A`$-module action $`ax:=\varphi (a)x`$.
We have found it convenient to signify right-Hilbert bimodule isomorphisms using diagrams: given right-Hilbert bimodules $`{}_{A}{}^{}X_{B}^{}`$, $`{}_{B}{}^{}Y_{C}^{}`$, and $`{}_{A}{}^{}Z_{C}^{}`$, when we say the diagram
commutes, we mean $`X_BYZ`$ as right-Hilbert $`AB`$ bimodules, and similarly for rectangular diagrams, etc.
If $`{}_{A}{}^{}X_{B}^{}`$ and $`{}_{C}{}^{}Y_{D}^{}`$ are right-Hilbert bimodules and $`\varphi :AC`$ and $`\psi :BD`$ are $`C^{}`$-homomorphisms, a linear map $`\mathrm{\Phi }:XY`$ is a *right-Hilbert bimodule homomorphism* with *coefficient maps* $`\varphi `$ and $`\psi `$ if it preserves the bimodule actions and the right inner product, that is,
* $`\mathrm{\Phi }(ax)=\varphi (a)\mathrm{\Phi }(x)`$,
* $`\mathrm{\Phi }(xb)=\mathrm{\Phi }(x)\psi (b)`$, and
* $`\mathrm{\Phi }(x),\mathrm{\Phi }(y)_D=\psi (x,y,_B)`$
for all $`aA`$, $`bB`$, and $`x,yX`$. If $`\varphi `$ and $`\psi `$ are isomorphisms then so is $`\mathrm{\Phi }`$, and if $`X`$ and $`Y`$ are imprimitivity bimodules then $`\mathrm{\Phi }`$ is an imprimitivity bimodule map.
An easy modification of \[16, Lemma 2.5\] shows that the $`B`$-linearity condition (ii) is redundant; in fact, if $`X_0`$ and $`A_0`$ are dense subspaces of $`X`$ and $`A`$, respectively, and if $`\mathrm{\Phi }:X_0Y`$ is a linear map satisfying (i) and (iii) above for all $`x,yX_0`$ and $`aA_0`$, then $`\mathrm{\Phi }`$ uniquely extends to a right-Hilbert bimodule homomorphism of $`{}_{A}{}^{}X_{B}^{}`$ into $`{}_{C}{}^{}Y_{D}^{}`$. Indeed, if $`SX`$ linearly spans $`X_0`$, $`TA`$ linearly spans $`A_0`$, $`\mathrm{\Phi }`$ satisfies (i) and (iii) on $`S`$ and $`T`$ and extends linearly to $`X_0`$, then the same conclusion holds. We will repeatedly use this fact without comment.
Given a right-Hilbert bimodule homomorphism $`{}_{A}{}^{}X_{B}^{}{}_{C}{}^{}Y_{D}^{}`$ with surjective coefficient maps $`\varphi :AC`$ and $`\psi :BD`$, by \[7, Lemma 5.3\] the diagram
commutes.
### Bimodule crossed products
Let $`Z`$ be an $`AB`$ imprimitivity bimodule, and let $`\eta :ZZC^{}(G)`$ be a bimodule coaction which is compatible with coactions $`\delta `$ on $`A`$ and $`ϵ`$ on $`B`$ (see ). Then the spectral subspaces $`\{Z_s:sG\}`$ densely span $`Z`$, and the bimodule crossed product $`Z\times _\eta G`$ is densely spanned by the pairs $`\{(x_s,t):s,tG,x_sZ_s\}`$, and is an $`A\times _\delta GB\times _ϵG`$ imprimitivity bimodule with operations given on the generators by
$`(a_r,s)(x_t,u)`$ $`=(a_rx_t,u)\text{if }s=tu\text{ (and }0\text{ if not)}`$
$`{}_{A\times G}{}^{}(x_r,s),(y_t,u)`$ $`=({}_{A}{}^{}x_r,y_t,tu)\text{if }s=u\text{ (and }0\text{ if not)}`$
$`(x_r,s)(b_t,u)`$ $`=(x_rb_t,u)\text{if }s=tu\text{ (and }0\text{ if not)}`$
$`(x_r,s),(y_t,u)_{B\times G}`$ $`=(x_r,y_t_B,u)\text{if }rs=tu\text{ (and }0\text{ if not)}.`$
The dual action $`\widehat{\eta }`$ of $`G`$ on $`Z\times _\eta G`$ is given on the generators by
$$\widehat{\eta }_t(x_r,s)=(x_r,st^1).$$
Similarly, if $`\gamma `$ is an action of $`G`$ on $`Z`$ which is compatible with actions $`\alpha `$ on $`A`$ and $`\beta `$ on $`B`$, then the crossed product $`Z\times _\gamma G`$ is densely spanned by the Cartesian product $`Z\times G`$, and is an $`A\times _\alpha GB\times _\beta G`$ imprimitivity bimodule with operations given on the generators by
$`(a,r)(x,s)`$ $`=(a\alpha _r(x),rs)`$
$`{}_{A\times G}{}^{}(x,r),(y,s)`$ $`=({}_{A}{}^{}x,\alpha _{rs^1}(y),rs^1)`$
$`(x,r)(b,s)`$ $`=(x\alpha _r(b),rs)`$
$`(x,r),(y,s)_{B\times G}`$ $`=(\alpha _{r^1}(x,y_B),r^1s).`$
The dual coaction $`\widehat{\gamma }`$ of $`G`$ on $`Z\times _\gamma G`$ is given on the generators by
$$\widehat{\gamma }(x,r)=(x,r)r.$$
### Imprimitivity theorems
Let $`(A,G,\delta )`$ be a maximal discrete coaction, and let $`N`$ be a normal subgroup of $`G`$. By the version of Mansfield’s imprimitivity theorem due to Echterhoff and the second author \[3, Theorem 3.1\], there exists an $`A\times _\delta G\times _{\widehat{\delta }|}NA\times _{\delta |}G/N`$ imprimitivity bimodule $`Y_{G/N}^G(A)`$. Mansfield’s bimodule is densely spanned by the Cartesian product $`𝒜\times G`$, with operations given on the generators by
$`(a_r,s,n)(b_t,u)`$ $`=(a_rb_t,un^1)\text{if }sn=tu\text{ (and }0\text{ if not)}`$
$`{}_{A\times G\times N}{}^{}(a_r,s),(b_t,u)`$ $`=(a_rb_t^{},ts,s^1u)\text{if }sN=uN\text{ (and }0\text{ if not)}`$
$`(a_r,s)(b_t,uN)`$ $`=(a_rb_t,t^1s)\text{if }sN=tuN\text{ (and }0\text{ if not)}`$
$`(a_r,s),(b_t,u)_{A\times G/N}`$ $`=(a_r^{}b_t,uN)\text{if }rs=tu\text{ (and }0\text{ if not)}`$
It is easy to see that $`Y(A)`$ is functorial in the sense that if $`(B,G,ϵ)`$ is another coaction and $`\varphi :AB`$ is an equivariant homomorphism then $`(a_s,t)(\varphi (a_s),t)`$ extends to an imprimitivity bimodule homomorphism $`Y(A)Y(B)`$ with coefficient homomorphisms $`\varphi \times G\times N`$ and $`\varphi \times G/N`$.
Dually, for an action $`(B,G,\alpha )`$ and a normal subgroup $`N`$ of $`G`$, Green’s imprimitivity theorem \[5, Proposition 3\] provides a $`B\times _\alpha G\times _{\widehat{\alpha }|}G/NB\times _{\alpha |}N`$ imprimitivity bimodule $`X_N^G(B)`$. Green’s bimodule is densely spanned by the Cartesian product $`B\times G`$, with operations given on the generators by
$`(a,r,sN)(b,t)`$ $`=(a\alpha _r(b),rt)\text{if }sN=tN\text{ (and }0\text{ if not)}`$
$`{}_{B\times G\times G/N}{}^{}(a,r),(b,s)`$ $`=(a\alpha _{rs^1}(b^{}),rs^1,sN)`$
$`(a,r)(b,n)`$ $`=(a\alpha _r(b),rn)`$
$`(a,r),(b,s)_{B\times N}`$ $`=(\alpha _{r^1}(a^{}b),r^1s)\text{if }rN=sN\text{ (and }0\text{ if not)}`$
$`X(B)`$ is functorial in the sense that if $`(C,G,\beta )`$ is another action and $`\varphi :BC`$ is an equivariant homomorphism then $`(b,s)(\varphi (b),s)`$ extends to an imprimitivity bimodule homomorphism $`X(B)X(C)`$ with coefficient homomorphisms $`\varphi \times G\times G/N`$ and $`\varphi \times N`$.
If the coaction $`(A,G,\delta )`$ is twisted over $`G/N`$, then the quotient map $`q:A\times _\delta GA\times _{\delta ,G/N}G`$ restricts to a surjection $`q|:A\times _{\delta |}G/NA`$, giving us an ideal $`\mathrm{ker}q|`$ of $`A\times G/N`$. Inducing across Mansfield’s bimodule $`Y`$ via the Rieffel correspondence gives an ideal $`Y\text{-}\mathrm{Ind}(\mathrm{ker}q|)`$ of $`A\times _\delta G\times _{\widehat{\delta }|}N`$, and \[12, Theorem 4.1\] shows that this ideal is precisely $`\mathrm{ker}(q\times N)`$. Rieffel’s theory thus gives an $`A\times _{\delta ,G/N}G\times _{\stackrel{~}{\delta }}NA`$ imprimitivity bimodule
$$Z_{G/N}^G(A):=Y/(Y\mathrm{ker}q|).$$
Moreover, the diagram
commutes.
Dually, if the action $`(B,G,\alpha )`$ is twisted over $`N`$, by \[5, Corollary 5\] we have a $`B\times _{\alpha ,N}G\times _{\widehat{\alpha }|}G/NB`$ imprimitivity bimodule
$$W_N^G(B):=X/(X\mathrm{ker}q|),$$
and we get a commutative diagram
If $`ϵ`$ is a maximal coaction of the quotient $`G/N`$ on $`A`$, the recent imprimitivity theorem for induced coactions \[2, Theorem 4.1\] gives an $`\mathrm{Ind}A\times _{\mathrm{Ind}ϵ}GA\times _ϵG/N`$ imprimitivity bimodule $`U=U(A)`$ densely spanned by the subset
$$\{(a_{sN},t):sNG/N,a_{sN}A_{sN},tG\}$$
of the Cartesian product $`A\times G`$, with operations given on the generators by
$`(a_{sN},t,r)(b_{uN},v)`$ $`=(a_{sN}b_{uN},tv)\text{if }r=v\text{ (and }0\text{ if not)}`$
$`(a_{sN},t)(b_{uN},vN)`$ $`=(a_{sN}b_{uN},t)\text{if }tN=suvN\text{ (and }0\text{ if not)}`$
$`{}_{\mathrm{Ind}A\times G}{}^{}(a_{sN},t),(b_{uN},v)`$ $`=(a_{sN}b_{uN}^{},tv^1,v)\text{if }s^1tN=u^1vN\text{ (and }0\text{ if not)}`$
$`(a_{sN},t),(b_{uN},v)_{A\times G/N}`$ $`=(a_{sN}^{}b_{uN},u^1vN)\text{if }t=v\text{ (and }0\text{ if not)}.`$
$`U(A)`$ is functorial in the sense that if $`(B,G/N,\eta )`$ is another coaction and $`\varphi :AB`$ is an equivariant homomorphism then $`(a_sN,t)(\varphi (a_{sN}),t)`$ extends to an imprimitivity bimodule homomorphism $`U(A)U(B)`$ with coefficient homomorphisms $`\mathrm{Ind}\varphi \times G`$ and $`\varphi \times G/N`$.
Dually, if $`\beta `$ is an action of the normal subgroup $`N`$ on $`B`$, the imprimitivity theorem for induced actions (sometimes attributed to Green \[5, Theorem 17\]) gives an $`\mathrm{Ind}B\times _{\mathrm{Ind}\beta }GB\times _\beta N`$ imprimitivity bimodule $`V=V(B)`$ densely spanned by the Cartesian product $`B\times G`$, with operations given on the generators by
$`([tr,b],t)(c,r)`$ $`=(bc,tr)`$
$`(b,s)(c,n)`$ $`=(\beta _{n^1}(bc),sn)`$
$`{}_{\mathrm{Ind}B\times G}{}^{}(b,s),(c,t)`$ $`=([s,bc^{}],st^1)`$
$`(b,s),(c,sh)_{B\times N}`$ $`=(b^{}\beta _h(c),h)`$
$`V(B)`$ is functorial in the sense that if $`(C,N,\beta )`$ is another action and $`\varphi :BC`$ is an equivariant homomorphism then $`(b,s)(\varphi (b),s)`$ extends to an imprimitivity bimodule homomorphism $`V(B)V(C)`$ with coefficient homomorphisms $`\mathrm{Ind}\varphi \times G`$ and $`\varphi \times N`$.
## 3. The Mansfield-Green Triangle
In this section we show a curious duality between the Mansfield and Green imprimitivity theorems. Theorem 3.1 will show that, roughly speaking, and modulo crossed product duality, Mansfield and Green induction are inverse processes.
Let $`(A,G,\delta )`$ be a maximal discrete coaction, and let $`N`$ be a normal subgroup of $`G`$. Not only do we have Mansfield’s $`A\times _\delta G\times _{\widehat{\delta }|}NA\times _{\delta |}G/N`$ imprimitivity bimodule $`Y_{G/N}^G(A)`$, but also, replacing $`N`$ by $`G`$, an $`A\times _\delta G\times _{\widehat{\delta }}GA`$ imprimitivity bimodule $`Y_{G/G}^G(A)`$. There is a $`\widehat{\widehat{\delta }}\delta `$ compatible coaction $`\delta _Y`$ of $`G`$ on $`Y_{G/G}^G(A)`$ (\[3, Remark 3.2\]) determined by
$$\delta _Y(a_r,s)=(a_r,s)s^1.$$
###### Theorem 3.1.
Let $`(A,G,\delta )`$ be a maximal coaction and let $`N`$ be a normal subgroup of $`G`$. Then the diagram
commutes.
###### Proof.
There is an imprimitivity bimodule isomorphism
$$\mathrm{\Phi }:Y(A)\times G/N\stackrel{~}{Y(A)}\stackrel{}{}X(A\times G)$$
defined on the generators by
$$\mathrm{\Phi }((a_r,s,tN)\stackrel{~}{(b_u,v)})=(a_rb_u^{},us,s^1v)\text{if }tN=vN\text{ (and }0\text{ if not)},$$
since straightforward calculations verify that the above mapping on generators preserves the left action and the right inner product. ∎
###### Remark 3.2.
To motivate the formula for $`\mathrm{\Phi }`$, note that
$$\mathrm{\Phi }((a_r,s,tN)\stackrel{~}{(c_u,v)})={}_{A\times G\times G}{}^{}(a_r,s),(c_u,v)\text{if }tN=vN\text{ (and }0\text{ if not)},$$
where $`(a_r,s)`$ and $`(c_u,v)`$ are viewed as elements of $`Y_{G/G}^G(A)`$, and the inner product is viewed as taking values in $`X_N^G(A\times _\delta G)`$.
The next two results will not be needed until Section 5; we include them here because they don’t involve twists, and are of general interest.
###### Proposition 3.3.
Let $`(A,G,\delta )`$ be a coaction, and let $`N`$ be a normal subgroup of $`G`$. Then
$$A\times _\delta G\times _{\widehat{\delta }}G\times _{\widehat{\widehat{\delta }}|}G/NA\times _{\delta |}G/N\times _{\delta ^{\mathrm{dec}}}G\times _{\widehat{\delta ^{\mathrm{dec}}}}G.$$
###### Proof.
Straightforward calculations verify that the mapping
$$(a_r,s,t,uN)(a_r,stuN,s,t)$$
on the generators preserves the operations, hence extends to a $`C^{}`$-isomorphism. ∎
###### Proposition 3.4.
Let $`(A,G,\delta )`$ be a coaction and $`N`$ a normal subgroup of $`G`$. Then the diagram
commutes, where the isomorphism is that of Proposition 3.3.
###### Proof.
There is an imprimitivity bimodule isomorphism
$$\mathrm{\Theta }:Y(A)\times G/N\stackrel{}{}Y(A\times G/N)$$
defined on the generators by
$$\mathrm{\Theta }(a_r,s,tN)=(a_r,tN,s),$$
since straightforward calculations verify that the above mapping on generators preserves the left action and the right inner product. ∎
## 4. The Dual Triangle
The results of this section are dual to those of the previous section, in the sense that actions correspond to coactions, Green bimodules correspond to Mansfield bimodules, and subgroups correspond to quotient groups. The only additional apparatus we need is to observe that if $`(B,G,\alpha )`$ is a discrete action then there is an $`\widehat{\widehat{\alpha }}\alpha `$ compatible action $`\alpha ^X`$ of $`G`$ on $`X_e^G(B)`$ given on the generators by
$$\alpha _r^X(a,s)=(a,sr^1).$$
###### Theorem 4.1.
Let $`(B,G,\alpha )`$ be an action an let $`N`$ be a normal subgroup of $`G`$. Then the diagram
commutes.
###### Proof.
There is an imprimitivity bimodule isomorphism
$$\mathrm{\Phi }:(X_e^G(B)\times N)\stackrel{~}{X_N^G}(B)\stackrel{}{}Y_{G/N}^G(B\times G)$$
defined on the generators by
$$\mathrm{\Phi }((a,r,n)\stackrel{~}{(b,s)})=(a\alpha _{rns^1}(b^{}),rns^1,sn^1),$$
since straightforward calculations verify that the above mapping on generators preserves the left action and the right inner product. ∎
###### Remark 4.2.
To motivate the formula for $`\mathrm{\Phi }`$, note that
$$\mathrm{\Phi }((a,r,n)\stackrel{~}{(b,s)})={}_{B\times G\times G}{}^{}(a,r),\alpha _n^X(b,s),$$
where $`(a,r)`$ and $`(b,s)`$ are viewed as elements of $`X_e^G(B)`$, and the inner product is viewed as taking values in $`Y_{G/N}^G(B\times _\alpha G)`$.
In analogy with the previous section, the next two results will not be needed until Section 6; they are presented here for convenience and general interest.
###### Proposition 4.3.
Let $`(B,G,\alpha )`$ be an action and let $`N`$ be a normal subgroup of $`G`$. Then
$$B\times _\alpha G\times _{\widehat{\alpha }}G\times _{\widehat{\widehat{\alpha }}|}NB\times _{\alpha |}N\times _{\alpha ^{\mathrm{dec}}}G\times _{\widehat{\alpha ^{\mathrm{dec}}}}G.$$
###### Proof.
Straightforward calculations verify that the mapping
$$(a,r,s,n)(a,rsns^1r^1,rsn^1s^1,sn)$$
on the generators preserves the operations, hence extends to a $`C^{}`$-isomorphism. ∎
###### Proposition 4.4.
Let $`(B,G,\alpha )`$ be an action and let $`N`$ be a normal subgroup of $`G`$. Then the diagram
commutes.
###### Proof.
There is an imprimitivity bimodule isomorphism
$$\mathrm{\Theta }:X(B)\times N\stackrel{}{}X(B\times N)$$
defined on the generators by
$$\mathrm{\Theta }(a,r,n)=(a,rnr^1,r),$$
since straightforward calculations verify that the above mapping on generators preserves the left action and the right inner product. ∎
## 5. The twisted Mansfield-Green square
Let $`(A,G,\delta )`$ be a maximal discrete coaction, and let $`N`$ be a normal subgroup of $`G`$. Combining Theorem 3.1 and Corollary 3.4, we get a commutative rectangle
(5.1)
Now suppose the coaction $`\delta `$ is twisted over $`G/N`$. Then the top arrow of the diagram (5.1) has
as a quotient. The imprimitivity bimodules in (5.1) above determine corresponding ideals of the bottom corners, and we can form a quotient commutative rectangle with upper right corner $`A`$. What happens to the rest of the diagram? We will answer this question in the present section.
However, we first modify the lower left corner of the diagram (5.1); the action $`\stackrel{~}{\delta }`$ of $`N`$ on $`A\times _{G/N}G`$ does not extend to $`G`$, so the Green bimodule $`X_N^G`$ on the left edge of (5.1) will not pass to a Green bimodule in the quotient. Rather, it will be more appropriate to use the bimodule arising from the imprimitivity theorem for induced actions. The action $`\stackrel{~}{\delta }`$ of $`N`$ on $`A\times _{G/N}G`$ induces to an action $`\mathrm{Ind}\stackrel{~}{\delta }`$ of $`G`$ on the induced algebra $`\mathrm{Ind}(A\times _{G/N}G)`$, and we have an $`\mathrm{Ind}(A\times _{G/N}G)\times _{\mathrm{Ind}\stackrel{~}{\delta }}G(A\times _{G/N}G)\times _{\stackrel{~}{\delta }}N`$ imprimitivity bimodule $`V_N^G(A\times _{G/N}G)`$.
It is shown in \[14, Theorem 4.4\] that
$$(a_sp_{tN},r)[r^1t,[a_s,t]]$$
extends to an isomorphism $`A\times _\delta G\mathrm{Ind}(A\times _{\delta ,G/N}G)`$ which is equivariant for the actions $`\widehat{\delta }`$ and $`\mathrm{Ind}\stackrel{~}{\delta }`$ of $`G`$. Turning this around and integrating up, we get an isomorphism
$$\mathrm{Ind}(A\times _{\delta ,G/N}G)\times _{\mathrm{Ind}\stackrel{~}{\delta }}GA\times _\delta G\times _{\widehat{\delta }}G.$$
###### Theorem 5.1.
If $`(A,G,\delta )`$ is a maximal discrete coaction which is twisted over $`G/N`$, the diagram
commutes.
###### Proof.
The desired diagram is the inner rectangle of the diagram
(5.2)
(Here and in Diagram (5.3) the action and coaction symbols have been omitted for clarity.) We will show how to fill in the bottom arrow so that each of the outer rectangle and the top, bottom, left, and right quadrilaterals commute. Since $`\mathrm{Ind}q\times G`$ is surjective, the result will then follow from standard bimodule techniques.
Consider the diagram
(5.3)
Since the action $`\widehat{\delta }|`$ of $`N`$ extends to $`G`$, it follows from standard facts concerning induced actions that
$$([s,a_t,r],u)(a_t,rs^1,u,u^1sN)$$
extends to an isomorphism $`\mathrm{Ind}(A\times _\delta G)\times _{\mathrm{Ind}\widehat{\delta }|}GA\times _\delta G\times _{\widehat{\delta }}G\times _{\widehat{\widehat{\delta }}|}G/N`$. Then an easy check on the generators shows that
$$(a_s,t,r)(a_s,tr^1,r)$$
extends to an isomorphism $`V(A\times _\delta G)X(A\times _\delta G)`$ of $`\mathrm{Ind}(A\times _\delta G)\times _{\mathrm{Ind}\widehat{\delta }|}GA\times _\delta G\times _{\widehat{\delta }|}N`$ imprimitivity bimodules. This shows the left triangle of the diagram (5.3) commutes. The inner quadrilateral in (5.3) is the commutative diagram (5.1). We *define* the isomorphism $`\mathrm{Ind}(A\times _\delta G)\times _{\mathrm{Ind}\widehat{\delta }|}GA\times _{\delta |}G/N\times _{\delta ^{\mathrm{dec}}}G\times _{\widehat{\delta ^{\mathrm{dec}}}}G`$ at the bottom arrow of (5.3) so that the bottom triangle commutes. On the generators, this isomorphism is given by
$$([s,a_t,r],u)(a_t,rN,rs^1,u).$$
Thus the outer rectangle in the diagram (5.2) commutes.
We noticed in Section 2 that the top quadrilateral in (5.2) commutes.
For the right quadrilateral, $`Y(B)`$ is functorial in $`B`$, so the homomorphism $`q|:A\times _{\delta |}G/NA`$ yields an imprimitivity bimodule homomorphism $`Y(q|):Y(A\times _{\delta |}G/N)Y(A)`$ with the desired coefficient homomorphisms. By \[7, Lemma 5.3\], this implies the quadrilateral commutes.
Similarly, the left quadrilateral commutes by functoriality of $`V(B)`$: the homomorphism $`q:A\times _\delta GA\times _{\delta ,G/N}G`$ yields an imprimitivity bimodule homomorphism $`V(q):V(A\times _\delta G)V(A\times _{\delta ,G/N}G)`$ with the desired coefficient homomorphisms.
Finally, the bottom quadrilateral in (5.2) commutes by a routine computation on the generators. ∎
## 6. The twisted dual square
In this section we introduce a twist into Theorem 4.1, just as in the preceding section we threw a twist into Theorem 3.1; unsurprisingly, the development will closely parallel that of Section 5.
Let $`(B,G,\alpha )`$ be a discrete action which is twisted over a normal subgroup $`N`$ in the sense of . Theorem 4.1 and Corollary 4.4 together give a commutative rectangle
(6.1)
As in the preceding section, in order to form a suitable quotient diagram we need to replace the lower left corner by an induced algebra.
The dual coaction $`\stackrel{~}{\alpha }`$ of $`G/N`$ on the twisted crossed product $`B\times _{\alpha ,N}G`$ induces to a coaction $`\mathrm{Ind}\stackrel{~}{\alpha }`$ of $`G`$ on the induced algebra $`\mathrm{Ind}(B\times _{\alpha ,N}G)`$, and we have an $`\mathrm{Ind}(B\times _{\alpha ,N}G)\times _{\mathrm{Ind}\stackrel{~}{\alpha }}G(B\times _{\alpha ,N}G)\times _{\stackrel{~}{\alpha }}G/N`$ imprimitivity bimodule $`U_{G/N}^G(B\times _{\alpha ,N}G)`$.
It is shown in \[2, Theorem 5.6\] that
$$(b,s)([b,s],s)$$
extends to an isomorphism $`B\times _\alpha G\mathrm{Ind}(B\times _{\alpha ,N}G)`$ which is equivariant for the coactions $`\widehat{\alpha }`$ and $`\mathrm{Ind}\stackrel{~}{\alpha }`$ of $`G`$. Turning this around and integrating up, we get an isomorphism
$$\mathrm{Ind}(B\times _{\alpha ,N}G)\times _{\mathrm{Ind}\stackrel{~}{\alpha }}GB\times _\alpha G\times _{\widehat{\alpha }}G.$$
###### Theorem 6.1.
If $`(B,G,\alpha )`$ is a discrete action which is twisted over $`N`$, the diagram
commutes.
###### Proof.
The desired diagram is the inner rectangle of the diagram
(6.2)
Consider the diagram
(6.3)
It follows from \[3, Remark 3.3\] that the map
$$(b,s,sn,t)(b,s,nt,t^1n^1t)$$
extends to an isomorphism $`\mathrm{Ind}(B\times _\alpha G)\times _{\mathrm{Ind}\widehat{\alpha }|}GB\times _\alpha G\times _{\widehat{\alpha }}G\times _{\widehat{\widehat{\alpha }}|}N`$, and this serves as the left-hand coefficient map for an isomorphism $`U(B\times _\alpha G)Y(B\times _\alpha G)`$, hence the left triangle of the diagram (6.3) commutes. The inner quadrilateral is the commutative diagram (6.1). We define the isomorphism $`\mathrm{Ind}(B\times _\alpha G)\times _{\mathrm{Ind}\widehat{\alpha }|}GB\times _{\alpha |}N\times _{\alpha ^{\mathrm{dec}}}G\times _{\widehat{\alpha ^{\mathrm{dec}}}}G`$ at the bottom arrow of (6.3) so that the bottom triangle commutes. On the generators, this isomorphism is given by
$$(b,s,t,r)(b,st^1,t,r).$$
Thus the outer rectangle in the diagram (6.2) commutes.
We noticed in Section 2 that the top quadrilateral in (6.2) commutes.
$`X(A)`$ is functorial in $`A`$, so the homomorphism $`q|:B\times _{\alpha |}NB`$ yields an imprimitivity bimodule homomorphism $`X(q|):X(B\times _{\alpha |}N)X(B)`$ with the desired coefficient homomorphisms, so the right quadrilateral commutes.
Similarly, the left quadrilateral commutes because by functoriality of $`U(A)`$ the homomorphism $`q:B\times _\alpha GB\times _{\alpha ,N}G`$ yields an imprimitivity bimodule homomorphism $`U(q):U(B\times _\alpha G)U(B\times _{\alpha ,N}G)`$ with the desired coefficient homomorphisms.
Finally, the bottom quadrilateral in (6.2) commutes by a routine computation on the generators. ∎
## 7. Ng’s Bimodule
We now return to the comparison between Ng’s bimodule and Mansfield’s, beginning with maximal coactions and full crossed products. In this context, by “Ng’s bimodule” we mean the bimodule gotten from the lower three sides of Diagram (7.1); the map $`\mathrm{\Theta }`$ will be defined in the proof of Theorem 7.1 by a construction parallel to Ng’s.
###### Theorem 7.1.
If $`(A,G,\delta )`$ is a maximal coaction of a discrete group $`G`$ and $`N`$ is a normal subgroup of $`G`$, then Ng’s bimodule is isomorphic to Mansfield’s; that is, the diagram
(7.1)
commutes.
###### Proof.
The desired diagram is the outer rectangle of
The upper left triangle commutes by Theorem 3.1, so we must show the lower right triangle commutes.
We construct the isomorphism $`\mathrm{\Theta }`$ as a composition
(7.2)
$$\begin{array}{c}A\times _\delta G\times _{\widehat{\delta }}G\times _{\widehat{\widehat{\delta }}|}G/N\stackrel{\mathrm{\Theta }_1}{}(A𝒦)\times _{ϵ_1}G/N\stackrel{\mathrm{\Theta }_2}{}(A𝒦)\times _{ϵ_2}G/N\hfill \\ \hfill \stackrel{\mathrm{\Theta }_3}{}(A\times _{\delta |}G/N)𝒦.\end{array}$$
Here, $`ϵ_2`$ is the coaction
$$ϵ_2:=(\mathrm{id}\sigma )(\delta |\mathrm{id})$$
of $`G/N`$ on $`A𝒦`$, where $`\sigma :AC^{}(G/N)𝒦A𝒦C^{}(G/N)`$ is the flip isomorphism. It follows from \[13, Lemma 1.16 (b)\] (see also ) that there is an isomorphism $`\mathrm{\Theta }_3`$ of $`(A𝒦)\times _{ϵ_2}G/N`$ onto $`(A\times _{\delta |}G/N)𝒦`$ defined on the generators by
$$\mathrm{\Theta }_3(a_{sN}b,tN)=(a_{sN},tN)b.$$
Perpetuating our perverse numbering scheme, we use $`ϵ_2`$ to define the coaction
$$ϵ_1:=\mathrm{Ad}(1U)ϵ_2$$
of $`G/N`$ on $`A𝒦`$, where $`U`$ is the unitary element of $`M(c_0(G)C^{}(G/N))M(𝒦C^{}(G/N))`$ determined by the bounded function $`U(s)=sN`$ from $`G`$ to $`C^{}(G/N)`$. It is easy to see that $`U`$ is an $`ϵ_1`$-cocycle (more precisely, the obvious analogue for full coactions of the more usual cocycles for reduced coactions—see ). It follows from \[9, Theorem 2.9\] (also see \[14, Proposition 2.8\]) that there is an isomorphism $`\mathrm{\Theta }_2`$ of $`(A𝒦)\times _{ϵ_1}G/N`$ onto $`(A𝒦)\times _{ϵ_2}G/N`$ defined by
$$\mathrm{\Pi }_2\mathrm{\Theta }_2=\mathrm{Ad}(\mathrm{id}_{A𝒦}\lambda )(U)\mathrm{\Pi }_1,$$
where $`\mathrm{\Pi }_i`$ is the regular representation of $`(A𝒦)\times _{ϵ_i}G/N`$ on $`\mathrm{}^2(G)\mathrm{}^2(G/N)`$ for $`i=1,2`$ (and $`A`$ is faithfully represented on a Hilbert space $``$).
Finally, from \[3, Equation (5.1) and Proposition 5.3\] we have the isomorphism $`\mathrm{\Phi }:A\times _\delta G\times _{\widehat{\delta }}GA𝒦`$ given on generators by
$$\mathrm{\Phi }(a_s,t,r)=a_s\lambda _sM_{\chi _t}\rho _r,$$
where $`\lambda `$ and $`\rho `$ are the left and right regular representations of $`G`$ and $`\chi _t`$ denotes the characteristic function of the singleton $`\{t\}`$. ($`\mathrm{\Phi }`$ is the isomorphism of Katayama’s duality theorem \[8, Theorem 8\], but for maximal coactions rather than reduced ones.) The arguments of , adapted to our context, show that $`\mathrm{\Phi }`$ is equivariant for the coactions $`\widehat{\widehat{\delta }}|`$ and $`ϵ_1`$ of $`G/N`$; we define $`\mathrm{\Theta }_1=\mathrm{\Phi }\times G/N`$ to be the corresponding isomorphism of the crossed products.
Careful study of the isomorphisms $`\mathrm{\Theta }_1`$, $`\mathrm{\Theta }_2`$, and $`\mathrm{\Theta }_3`$ now shows that the composition $`\mathrm{\Theta }:A\times _\delta G\times _{\widehat{\delta }}G\times _{\widehat{\widehat{\delta }}|}G/N\stackrel{}{}(A\times _{\delta |}G/N)𝒦`$ is given on the generators by
$$\mathrm{\Theta }(a_s,t,r,gN)=(a_s,trqN)\lambda _sM_{\chi _t}\rho _r.$$
Using this, straightforward calculations show that there is an $`A\times _\delta G\times _{\widehat{\delta }}G\times _{\widehat{\widehat{\delta }}|}G/NA\times _{\delta |}G/N`$ imprimitivity bimodule isomorphism
$$\mathrm{{\rm Y}}:Y_{G/G}^G(A)\times _{\delta _Y}G/N\stackrel{}{}(A\times _{\delta |}G/N)\mathrm{}^2(G)$$
defined on the generators by
$$\mathrm{{\rm Y}}(a_s,t,rN)=(a_s,rN)\chi _{st}.$$
###### Remark 7.2.
Taking $`N=G`$ in Theorem 7.1 shows that Katayama’s bimodule (by which we mean the bottom and right-hand sides of that rectangle, taken together) is isomorphic to Mansfield’s in this special case. This justifies the idea that Mansfield’s theorem “reduces to Katayama’s” when $`N=G`$, a fact which is well-known to the cognoscenti, but to our knowledge has not explicitly appeared in the literature.
To complete the connection with Ng’s theorem, we need to pass to *reduced* coactions and *amenable* subgroups in Diagram 7.1. Here we use $`Y(A,\delta )`$ to denote the $`A\times _\delta G\times _{\widehat{\delta }|}N`$$`A\times _{\delta |}G/N`$ imprimitivity bimodule provided by the original form of Mansfield’s Imprimitivity Theorem \[10, Theorem 27\]. The isomorphism $`\mathrm{\Theta }_r:A\times _\delta G\times _{\widehat{\delta }}G\times _{\widehat{\widehat{\delta }}|}G/N(A\times _{\delta |}G/N)𝒦(\mathrm{}^2)`$ is constructed as in Equation (7.2).
###### Corollary 7.3.
If $`(A,G,\delta )`$ is a reduced coaction of a discrete group $`G`$ and $`N`$ is an amenable normal subgroup of $`G`$, then Ng’s bimodule is isomorphic to Mansfield’s; that is, the diagram
commutes.
###### Proof.
Since $`G`$ is discrete, the coaction $`\delta `$ is automatically nondegenerate, so by \[13, Theorem 4.7\] there is a unique full coaction $`\delta ^f`$ of $`G`$ on $`A`$ whose reduction coincides with $`\delta `$, and then \[13, Proposition 3.8\] gives an isomorphism $`A\times _\delta G\stackrel{}{}A\times _{\delta ^f}G`$; it is easy to see that this isomorphism is equivariant for the dual actions. Then \[3, Proposition 5.3\] applies, giving a maximal coaction $`(A^m,G,\delta ^m)`$ (the “maximalization” of $`\delta ^f`$) and an equivariant surjection $`\mathrm{\Psi }:A^mA`$ whose integrated form $`\mathrm{\Psi }\times G:A^m\times _{\delta ^m}GA\times _{\delta ^f}G`$ is an isomorphism which is equivariant for the dual actions. Then $`\mathrm{\Psi }`$ is also equivariant for the restricted coactions $`\delta ^m|`$ and $`\delta ^f|`$, hence certainly gives a surjection
$$\mathrm{\Psi }\times G/N:A^m\times _{\delta ^m|}G/NA\times _{\delta ^f|}G/N.$$
Since $`\delta ^f`$ is the normalization of $`\delta ^m`$, \[3, Theorem 3.4\] tells us that, if $`I=\mathrm{ker}\mathrm{\Psi }\times G/N`$, then the ideal of $`A^m\times _{\delta ^m}G\times _{\widehat{\delta ^m}|}N`$ induced from $`I`$ via the Mansfield imprimitivity bimodule $`Y(A^m)`$ coincides with the kernel of the regular representation
$$A^m\times _{\delta ^m}G\times _{\widehat{\delta ^m}|}NA^m\times _{\delta ^m}G\times _{\widehat{\delta ^m}|,r}N,$$
so that $`Y/(YI)`$ is canonically an $`A^m\times _{\delta ^m}G\times _{\widehat{\delta ^m}|,r}N`$$`A\times _{\delta ^f|}G/N`$ imprimitivity bimodule. But $`N`$ is amenable, so the regular representation of $`A^m\times _{\delta ^m}G\times _{\widehat{\delta ^m}|}N`$ is faithful. Hence we must have $`I=\{0\}`$, so $`\mathrm{\Psi }\times G/N`$ is actually an isomorphism of $`A^m\times _{\delta ^m|}G/N`$ onto $`A\times _{\delta ^f|}G/N`$.
It is now clear from the constructions that the identity map on the ordered pairs $`\{(a_s,t):s,tG\}`$ extends to an isomorphism
$$Y(A^m)\stackrel{}{}Y(A,\delta ^f)$$
of the Mansfield imprimitivity bimodule $`Y(A^m)`$ used in the present paper onto the version of the Mansfield bimodule associated to the normal coaction $`\delta ^f`$ in . Thus the diagram
commutes. On the other hand, one of the main points of is that the diagram
commutes; combining these shows that the top quadrilateral of the diagram
commutes. Since the outer rectangle commutes by Theorem 7.1, and the left, right, and bottom quadrilaterals are easily seen to commute, we conclude that the inner rectangle commutes as well. ∎ |
warning/0002/nlin0002014.html | ar5iv | text | # Breaking time reversal symmetry in chaotic driven Rydberg atoms
## I INTRODUCTION
Sixteen years ago, Bohigas, Giannoni and Schmit, in a seminal paper , formulated the conjecture that quantum systems which are chaotic in the classical limit, generically have statistical properties of energy levels described by Random Matrix Theory (RMT) . While RMT had proven useful in characterizing nuclear spectra , the conjecture was quite surprising since it claimed an applicability of RMT statistics to deterministic systems with very few degrees of freedom. While some counterexamples exist (see e.g. ), the conjecture remains widely accepted and played a very stimulating role in quantum chaos studies (see for reviews ).
Depending on the symmetries of a given strongly chaotic system, the statistical properties of its quantum spectrum fall into one of the three classes known from RMT: orthogonal, unitary and symplectic. The orthogonal class is associated with Hamiltonians invariant with respect to some anti-unitary symmetry. Typically this is true for time-reversal invariant systems, but also for more complicated anti-unitary symmetry (generalized time-reversal symmetry) as the product of a time-reversal with some discrete reflection. Roughly speaking, this happens if a basis can be chosen where the matrix elements of the Hamiltonian are all real. The corresponding statistical ensemble of (real symmetric) random matrices is referred to as the Gaussian Orthogonal Ensemble (GOE). In the absence of any anti-unitary symmetry, the corresponding class of complex hermitian matrices is known as the Gaussian Unitary Ensemble (GUE). Finally, specific considerations apply to half-integer spin systems with (generalized) time-reversal symmetry. Indeed, the energy levels of these systems are systematically two-fold degenerate (this is the well-known Kramers degeneracy). In the presence of a geometrical symmetry like azimuthal symmetry, the Kramers degeneracy is hidden and the GOE statistics apply. With no geometrical symmetries, the Gaussian Symplectic Ensemble (GSE) has to be used. All these three classes of systems are characterized by level repulsion, that is zero probability of observing accidentally degenerate energy levels. This is to be contrasted with a generic behaviour of multidimensional integrable systems where close lying levels are uncorrelated and obey a Poisson statistics . For a generic Hamiltonian system, where typically chaotic and regular motions coexist, the statistical properties of energy levels are not universal and are expected to be intermediate between the two limiting cases, Poisson and RMT.
While the conjecture has been tested on a number of theoretical models (we refer the reader to reviews rather than numerous original papers), experimental verifications are much less abundant and restricted fully, as far as we know, to the orthogonal universality class . This is probably due to the fact that most experimental results are obtained for atomic and molecular spectra. In the presence of a uniform magnetic field, the time-reversal symmetry is broken, but a anti-unitary symmetry persists. Thus, observing GUE statistics should require well controlled inhomogeneous fields on the atomic scale, which has not been achieved.
An alternative is to study the eigenmodes not of the Schrödinger equation, but of a different although similar wave equation. The best examples are two-dimensional microwave billiards (in fact three-dimensional flat cavities below the cut-off frequency) where the classical Helmholtz wave equation is in fact equivalent to the time-independent Schrödinger equation. There, breaking of the time reversal invariance is possible by using magnetic devices and GUE-type statistics have been observed experimentally . Still it is desirable to have at hand a quantum system where manifestations of the generalized time reversal symmetry breaking is experimentally accessible. The aim of this paper is to discuss an example of such a system. As mentioned above, breaking all anti-unitary symmetries using static fields is not possible. Hence, the idea is to use a time-dependent field acting on an atom. With an oscillatory field alone, the product of time reversal by symmetry with respect to one of the polarization axis is an anti-unitary symmetry. Hence, GUE statistics requires to combine a microwave field and a static field.
The model we use, the Rydberg states of the hydrogen atom driven by a microwave field, has a long history of its own. The system attracted attention when it turned out that experimental results on the ionization probability as a function of the microwave field and frequency could be explained in terms of the underlying classical chaotic dynamics which results in a diffusive energy gain of the Rydberg electron and eventually to ionization. Many interesting phenomena (e.g. quantum localization) have been studied on this model both theoretically and experimentally (see e.g. review papers ). The first experiments involved linearly polarized microwaves allowing for a rough description of the ionization thresholds via a simple one-dimensional time-dependent model where the motion of the electron is restricted to the microwave field axis. For other polarizations (available in recent experiments ), such a simple model is not possible: the simplest dynamics may be two-dimensional where the electronic motion is restricted to the polarization plane. Such models have been investigated both for circular polarization (CP) and general elliptical polarization (EP) .
When not only the ionization thresholds, but also the more subtle details of the dynamics, have to be studied, a full three dimensional atom must be considered. For LP microwave, the conservation of the angular momentum projection onto the polarization axis, $`L_z`$, makes the dynamics effectively two-dimensional and time dependent. For CP case, while $`L_z`$ is not conserved, the transformation to the frame rotating with the microwave frequency removes the explicit oscillatory time-dependence, leading to a three-dimensional time independent problem. Both these simplifications are no longer possible in the general EP microwave field and the problem becomes truly three dimensional and time dependent, providing new challenges to the theory. It is this situation which we shall consider later on.
However, instead of discussing typical “large” microwave amplitude for fields which lead to an efficient excitation of Rydberg atoms to higher excited states and finally to ionization, we shall consider a “weak” perturbation of an atom where the exchange of energy between the field and the atom is negligible. The fact that the Coulomb problem is highly degenerate provides us with the possibility, as shown below, to reach chaotic dynamics even for infinitesimally weak external fields. This is not in contradiction with the Kolmogorov-Arnold-Moser (KAM) theorem since the latter precisely does not apply for degenerate situations. Those weak fields do not really excite the atom, but rather couple and mix among each other the $`n_0^2`$ states of a degenerate hydrogenic manifold with a given principal quantum number $`n_0`$. The degeneracy is lifted, states undergo small (but measurable) shifts that reflect the dynamics of the $`n_0`$ manifold. Such a situation is referred to as intramanifold dynamics (by contrast to a typical intermanifold coupling leading to atomic excitation and ionization). An intramanifold chaotic behaviour is quite attractive from the theoretical point of view. It yields an effective quantum Hamiltonian acting in the $`n_0`$ space, whose eigenvalues are the energy shifts, and is represented by a finite dimensional matrix of size $`n_0^2`$ (modulo the remaining point symmetries) with no cutoff errors. The semiclassical limit is realized by letting $`n_0\mathrm{}`$.
The first attempt to produce some intramanifold chaos used the hydrogen atom in uniform crossed magnetic and electric fields , which was later extended to arbitrary mutual orientations of the two fields . The authors, exploring second order perturbation theory, have observed signatures of intramanifold chaos in the quantum spectrum. However, the situation is somewhat complex: the first order term (in the two external fields) is always integrable; it is only when combined with a second order term of a comparable magnitude that some noticeable region of chaos can be created. But this implies the application of large fields, and higher order terms in the perturbative expansion are of importance. That, unfortunately, has been neglected by the authors. Moreover, the high electric field values used lead then to an extremely fast field ionization which blurs any long-time effects in the dynamics. Especially, there are no more bound states in these conditions, but only broad overlapping resonances, which implies that the statistical properties cannot be described by one of the three standard ensembles of random matrices .
Recently, we have proposed another system revealing chaotic intramanifold dynamics which does not suffer from the mentioned deficiency. This is an hydrogen atom driven resonantly (i.e. with the frequency $`\omega `$ which is an integer multiple of the Kepler frequency $`\omega _\mathrm{K}=1/n_0^3`$) by an elliptically polarized microwave with the possible addition of some static weak electric field. The latter allows us to break any generalized time-reversal symmetry of the system.
Since the driving field is periodic, by applying the Floquet theorem, one can find the eigenstates of the system (the so-called Floquet or dressed states). The eigenenergies of the system are referred to as quasienergies of the system and are defined modulo $`\mathrm{}\omega `$. As we have discussed shortly for a $`2:1`$ resonance , the statistical properties of the quasienergies reveal convincingly the symmetry breaking. Here we discuss the same system, i.e. an hydrogen atom driven by a resonant EP microwave field, both in the presence and in the absence of the static field.
The paper is organized as follows. In Sec. II we derive the quantum perturbative Hamiltonian for our system and its semiclassical counterpart. They are used in Sec. III to analyze the behavior of the atom in a pure microwave field and in Sec. IV for a microwave field combined with a static electric field. The summary and the future perspectives form the content of the concluding section.
## II PERTURBATION APPROACH
We consider a realistic three-dimensional hydrogen atom placed in a static electric field and driven by an elliptically polarized microwave field. We define the $`z`$-axis as perpendicular to the plane of the polarization of the microwave field. The Hamiltonian of the system, in atomic units, neglecting relativistic effects, assuming infinite mass of the nucleus, and employing dipole approximation reads:
$$H=\frac{𝐩^2}{2}\frac{1}{r}+F(x\mathrm{cos}\omega t+\alpha y\mathrm{sin}\omega t)+𝐄𝐫,$$
(1)
where $`F`$, $`\alpha `$ and $`\omega `$ stand, respectively, for the amplitude, degree of elliptical polarization and frequency of the microwave field while $`𝐄`$ denotes the static electric field.
As mentioned in the Introduction, we are interested in effects due to weak external fields. In the absence of any external field, the energy levels of the system are $`1/2n_0^2`$ with $`n_0`$ the principal quantum number; the degeneracy of the $`n_0`$ hydrogenic manifold is $`n_0^2.`$ Even very weak fields mix strongly states within the manifold; on the other hand, it is perfectly justified to treat coupling to other manifolds perturbatively. We shall do so first quantum-mechanically with the help of the effective Hamiltonian approach . Then, we construct its (semi-)classical equivalent. The latter allows us to study the character of the classical motion and search for parameters corresponding to chaotic dynamics.
### A Quantum perturbation method
For any time-periodic Hamiltonian, the Floquet theorem states that the most general solution of the Schrödinger equation can be written as a linear combination of the “Floquet states”, which are time-periodic functions and eigenstates of the so-called Floquet Hamiltonian $`_{\mathrm{F}loquet}`$ of the system:
$$_{\mathrm{F}loquet}\varphi (t)=\left(Hi\frac{}{t}\right)\varphi (t)=ϵ\varphi (t)$$
(2)
where $`ϵ`$ and $`\varphi (t)`$ are respectively the quasienergy and the Floquet state. The solutions of Eq. (2) have to satisfy the usual boundary conditions in configuration space and periodic boundary condition in the time coordinate (by construction, Floquet states are time-periodic). Hence, the Floquet Hamiltonian $`_{\mathrm{F}loquet}`$ acts on an extended Hilbert space containing also the time coordinate.
For studying the quasi-energy spectrum, is it convenient to choose, as a basis of the atomic Hilbert space, the Sturmian functions (see e.g. for details of this application) $`n,L,M^{(\mathrm{\Lambda })}`$, where $`L`$ and $`M`$ are the usual angular and magnetic quantum numbers respectively while $`nL+1`$ labels the radial functions whose number of nodes is $`nL1.`$ $`\mathrm{\Lambda }`$ is a scaling parameter (unit of length in configuration space) for the Sturmian functions. For $`\mathrm{\Lambda }=n_0,`$ the Sturmian functions with $`n=n_0`$ are the exact hydrogenic states, eigenstates of the unperturbed atom. As we intend to describe the dynamics inside the $`n_0`$-manifold of the hydrogen atom, we choose to keep $`\mathrm{\Lambda }=n_0`$ in all calculations.
Along the time coordinate, we choose the usual oscillating exponential functions as a basis of time-periodic functions. They are labeled with an integer index $`K`$ and defined by:
$$t|K=e^{iK\omega t}$$
(3)
The whole Hilbert space of the atom + periodic perturbation system (Floquet Hamiltonian) is spanned by the tensor product of the configuration space and time basis:
$$n,L,M,K=n,L,M^{(n_0)}|K.$$
(4)
In the dressed atom language, $`K`$ may be loosely identified with the number of photons exchanged between the atom and the field.
The Sturmian basis is not orthogonal, but it satisfies the following relation:
$$\underset{n,L,M,K}{}|n,L,M,Kn,L,M,K|\frac{1}{2r}=1$$
(5)
The advantage is that, when written in the Sturmian basis, all the matrix elements representing the various parts of the Floquet Hamiltonian have some strong selection rules. The selection rules on $`K`$ trivially come from the Fourier expansion of the time dependences; the selection rules on $`L`$ and $`M`$ originate from the angular dependence of the various operators. The selection rules on $`n`$ are far from obvious and are at the heart of the definition and properties of the Sturmian functions .
If we define:
$`H_0`$ $`=`$ $`{\displaystyle \frac{𝐩^2}{2}}{\displaystyle \frac{1}{r}}i{\displaystyle \frac{}{t}},`$ (6)
$`U`$ $`=`$ $`F(x\mathrm{cos}\omega t+\alpha y\mathrm{sin}\omega t),`$ (7)
$`V`$ $`=`$ $`𝐄𝐫,`$ (8)
we obtain the following selection rules for the matrices representing these operators in the Sturmian basis:
$$\begin{array}{ccccc}\mathrm{\Delta }n=0,\pm 1,\hfill & \mathrm{\Delta }L=0,\hfill & \mathrm{\Delta }M=0,\hfill & \mathrm{\Delta }K=0,\hfill & \text{for }H_0\text{ and unity operators},\hfill \\ \mathrm{\Delta }n=0,\pm 1,\pm 2\hfill & \mathrm{\Delta }L=\pm 1,\hfill & \mathrm{\Delta }M=0,\pm 1\hfill & \mathrm{\Delta }K=0,\hfill & \text{for }V,\hfill \\ \mathrm{\Delta }n=0,\pm 1,\pm 2,\hfill & \mathrm{\Delta }L=\pm 1,\hfill & \mathrm{\Delta }M=\pm 1,\hfill & \mathrm{\Delta }K=\pm 1,\hfill & \text{for }U.\hfill \end{array}$$
(9)
In addition, all matrix elements are known in closed forms and involve only square roots of rational functions of the various quantum numbers. Note that, because of the non orthogonal character of the Sturmian basis, calculating the Floquet quasi-energies requires to solve a generalized eigenvalue problem rather than a standard one. This is the price to pay for getting sparse matrices.
The exact calculation of the quasi-energies is possible only numerically. However, we are interested in the situation where both the static and the microwave fields are weak. Thus, a perturbative expansion is convenient. We will now perform it at the lowest non-vanishing order for each external field. Because, the zeroth order eigenstates are highly degenerate, we have to use degenerate perturbation theory. A convenient formulation is to use an effective Hamiltonian which, at any order of the calculation, has the same spectrum as the initial Hamiltonian, but acts only inside the degenerate hydrogenic manifold. The details of the method are given in . At first order, the calculation is trivial and the effective Hamiltonian $`H^{(1)}`$ is just the projection of the perturbation onto the manifold we are considering. If $`P`$ denotes the projector onto the degenerate $`(n_0,K=0)`$ manifold, it is simply:
$$H^{(1)}=P(U+V)P=PVP$$
(10)
since $`U`$ always changes $`K`$ by one unit. The non-zero matrix elements of $`H^{(1)}`$ are those of $`V`$ with $`\mathrm{\Delta }n=0.`$ Thus $`H^{(1)}`$ is proportional to the static electric field.
The lowest non-vanishing contribution of the microwave field is at second order. It has the following well-known formal expression:
$$H^{(2)}=PUQ\frac{1}{E_0H_0}QUP$$
(11)
where $`E_0=1/2n_0^2`$ is the unperturbed energy of the hydrogenic manifold and $`Q=1P`$ is the projector onto the subspace complementary to the hydrogenic manifold.
Explicit calculation of $`H^{(2)}`$ is not straightforward. Indeed, if one expands the $`1/(E_0H_0)`$ onto the eigenstates of $`H_0,`$ one obtains a infinite sum over the discrete states and continuum of the atomic spectrum. The trick is to use the non-orthogonal Sturmian basis defined above. Indeed, the projector $`P`$ has the following simple expression is this basis:
$$P=\underset{L,M}{}|n_0,L,M,0n_0,L,M,0|\frac{1}{2r},$$
(12)
and consequently
$$Q=\underset{L,M,(n,K)(n_0,0)}{}|n,L,M,Kn,L,M,K|\frac{1}{2r}.$$
(13)
The last step is to calculate the matrix element of the $`1/(E_0H_0)`$ operator in the Sturmian basis. It is simply accomplished by noting that the operator $`(E_0H_0)`$ is diagonal in $`L`$,$`M`$ and $`K`$ and tridiagonal in $`n`$ (i.e. connects only state $`n`$ to states $`n1`$,$`n`$ and $`n+1`$). Thus the matrix elements of $`1/(E_0H_0)`$ are simply obtained by solving a triadiagonal set of coupled equations in each $`(L,M,K)`$ subspace coupled to the initial state.
Finally the whole effective quantum Hamiltonian inside the $`n_0`$-manifold reads
$$H_{\mathrm{e}ff}=H^{(1)}+H^{(2)}.$$
(14)
This Hamiltonian takes into account the direct coupling between the levels due to the presence of the static electric field (the term proportional to $`E`$) and the indirect coupling through all levels of other manifolds, i.e. process of absorption and emission of microwave photons (the term proportional to $`F^2`$) . Because of the selection rules on the $`U`$ and $`V`$ operators and the simple algebraic structure of the effective Hamiltonian, $`H_{\mathrm{e}ff}`$ itself has the following selections rules: $`\mathrm{\Delta }n=0,\mathrm{\Delta }K=0`$ (by construction) and $`\mathrm{\Delta }L=0,\pm 1,\pm 2,`$ $`\mathrm{\Delta }M=0,\pm 1,\pm 2.`$ The diagonalization of $`H_{\mathrm{e}ff}`$, i.e. of the sparse banded matrix of dimension $`n_0^2`$, by standard routines, yields quasienergies of the system.
The method has been tested in limiting situations, e.g., for parallel weak static and linearly polarized fields where quasienergies resulting from the effective Hamiltonian could be compared with exact diagonalization values. The results for EP presented below are obtained for field amplitudes for which excellent agreement between exact and perturbative results have been found in the limiting cases.
### B Semiclassical perturbation method
The general prescription for the semiclassical quantization of a time-periodic system has been described in . Recently, we used a similar procedure for a LP microwave , for a static electric field parallel to a LP microwave and for a two-dimensional model atom in the EP microwave case .
The method requires first to define a classical Hamilton function for which one can use the usual semiclassical quantization rules. This is done by passing to the extended phase space, defining the momentum $`p_t`$ conjugate to the $`t`$ (time) variable. It yields the new classical Hamiltonian, $`_{\mathrm{F}loquet}=H+p_t`$ , being the classical analog of the quantum Floquet Hamiltonian, Eq. (2).
As the next step, we express the Hamiltonian in action-angle variables of the unperturbed Coulomb problem . Due to its high symmetry, several choices are possible. The standard solution is to consider the canonically conjugate pairs $`(J,\mathrm{\Theta })`$, $`(L,\mathrm{\Psi })`$ and $`(M,\mathrm{\Phi })`$. $`J`$ is the principal action (corresponding to the principal quantum number), i.e. the total action along an unperturbed Kepler elliptical trajectory of the electron. It is simply related to the size of the ellipse. The corresponding angle, $`\mathrm{\Theta }`$, determines the position of the electron on its elliptical trajectory and depends linearly on time, $`\mathrm{\Theta }\omega _Kt`$, for an unperturbed atom. $`L`$ is the angular momentum, $`\mathrm{\Psi }`$ the angle of rotation around the axis defined by the angular momentum vector. Similarly, $`M,\mathrm{\Phi }`$ denote the projection of the angular momentum on the laboratory $`z`$-axis and the angle of rotation around that axis, respectively. The shape of the ellipse is best described by its eccentricity $`e=\sqrt{1L^2/J^2}`$ while its orientation in the configuration space is determined by the Euler angles $`(\mathrm{\Phi },\mathrm{arccos}M/L,\mathrm{\Psi })`$ as defined by Goldstein .
Using these canonical coordinates, it is possible to write down the full Floquet Hamiltonian. We now specialize to the resonant case where the microwave frequency is an integer multiple of the Kepler frequency of the unperturbed electron, i.e. $`\omega _0=\omega /\omega _K=m`$. The corresponding action is:
$$J=n_0=\omega _K^{1/3}=\left(\frac{\omega }{m}\right)^{1/3}.$$
(15)
$`n_0`$ is interpreted as the principal quantum number of the quantum hydrogenic manifold where the resonance takes place.
In the absence of any external field, the variables $`(J,L,\mathrm{\Psi },M,\mathrm{\Phi })`$ are all constant while $`\mathrm{\Theta }`$ evolves linearly in time (see above). Hence, in the presence of weak external fields, the motion along the $`\mathrm{\Theta }`$ variable will be much faster than along the other coordinates and the secular perturbation theory can be used: it averages over the nonresonant terms and yields the approximate resonant Hamiltonian of the form
$$_{\mathrm{s}ec}=\frac{1}{2n_0^2}\frac{3m^2}{2n_0^4}\widehat{J}^2+F\mathrm{\Gamma }_m\mathrm{cos}(\widehat{\mathrm{\Theta }}\delta )+E\gamma +\widehat{p}_t$$
(16)
where
$`\gamma `$ $`=`$ $`{\displaystyle \frac{3}{2}}n_0^2[\mathrm{cos}\phi \mathrm{sin}\theta (\mathrm{cos}\mathrm{\Phi }\mathrm{cos}\mathrm{\Psi }{\displaystyle \frac{M}{L}}\mathrm{sin}\mathrm{\Phi }\mathrm{sin}\mathrm{\Psi })`$ (19)
$`+\mathrm{sin}\phi \mathrm{sin}\theta \left(\mathrm{sin}\mathrm{\Phi }\mathrm{cos}\mathrm{\Psi }+{\displaystyle \frac{M}{L}}\mathrm{cos}\mathrm{\Phi }\mathrm{sin}\mathrm{\Psi }\right)`$
$`+\mathrm{cos}\theta \sqrt{1{\displaystyle \frac{M^2}{L^2}}}\mathrm{sin}\mathrm{\Psi }]\sqrt{1{\displaystyle \frac{L^2}{n_0^2}}}.`$
and
$$\widehat{\mathrm{\Theta }}=m\mathrm{\Theta }\omega t,\widehat{J}=\frac{Jn_0}{m},\widehat{p}_t=p_t+\frac{\omega J}{m},$$
(20)
The secular variables $`\widehat{\mathrm{\Theta }},\widehat{J},\widehat{p_t}`$ are slowly varying variables obtained by substracting the unperturbed resonant quantities. $`\widehat{\mathrm{\Theta }}`$ represents the phase drift of the electron along the elliptical orbit and $`\widehat{J}`$ the distance (in action) to the exact resonance.
The orientation of the static field with respect to the $`z`$-axis is determined by the usual spherical angles, $`\phi `$ and $`\theta `$. The expression for $`\gamma `$ looks complicated, but it is actually nothing but the component of the average atomic dipole on the static field axis. Similarly, $`\mathrm{\Gamma }_m(L,\mathrm{\Psi },M,\mathrm{\Phi };\alpha )`$ and $`\delta (L,\mathrm{\Psi },M,\mathrm{\Phi };\alpha )`$ just represent the amplitude and the phase of the atomic dipole at the microwave frequency. They can be obtained simply from the Fourier components of the electron position on an unperturbed elliptical trajectory. The explicit, rather complicated formulae for $`\delta (L,\mathrm{\Psi },M,\mathrm{\Phi };\alpha )`$ and $`\mathrm{\Gamma }_m(L,\mathrm{\Psi },M,\mathrm{\Phi };\alpha )`$ are given by Eq. (2.15) and Eq. (2.16) of , respectively and are reproduced in Appendix A for the convenience of the reader.
The last stage is to quantize the system using the approximate Hamiltonian, Eq. (16). As any explicit time dependence has disappeared in the secular Hamiltonian, the quantization of $`\widehat{p}_t`$ is trivial. Taking into account that eigenstates of the Floquet Hamiltonian have to be time-periodic, this yields $`\widehat{p}_t=k\omega `$ (where $`k`$ is an integer number) which ensures the periodicity of the quasienergy spectrum with a period $`\omega `$.
The radial motion, i.e. in the $`(\widehat{J},\widehat{\mathrm{\Theta }})`$ effectively decouples from the secular motion of the elliptical trajectory, i.e. in the $`(L,\mathrm{\Psi },M,\mathrm{\Phi })`$ space . While considering the radial motion, the effective hamiltonian for the secular part is approximately constant (for a detailed discussion as well as possible counterexamples in some cases see ). In effect the quantization resembles in spirit the Born-Oppenheimer approximation. One may first quantize the radial motion keeping the secular motion frozen. The radial motion exhibits a pendulum-like dynamics whose quantum eigenvalues are given by the solutions of the Mathieu equation (see the similar treatment for one-dimensional systems ). Hence, we can define an effective Hamiltonian acting in a reduced $`(L,\psi ,M,\varphi )`$ phase space just replacing the $`(\widehat{J},\widehat{\mathrm{\Theta }})`$ part of $`_{\mathrm{s}ec}`$ by the quantized energy levels of the pendulum. In this paper, we are interested in the ground state of the pendulum, thus, the quantization of the radial motion yields the following effective Hamiltonian for the secular motion:
$$_{\mathrm{e}ff}=\frac{3m^2}{8n_0^4}a_0(q)+E\gamma ,$$
(21)
where the constant terms $`1/2n_0^2`$ and $`k\omega `$ are omitted ($`_{\mathrm{e}ff}`$ denotes the shift from the unperturbed energy level of the atom), and:
$$q=\frac{4n_0^4F}{3m^2}\mathrm{\Gamma }_m.$$
(22)
is a dimensionless parameter. $`a_0(q)`$ is the Mathieu parameter corresponding to the ground state of the pendulum. We can introduce scaled quantities:
$`F_0=n_0^4F`$ (23)
$`E_0=n_0^4E`$ (24)
$`L_0={\displaystyle \frac{L}{n_0}}`$ (25)
$`M_0={\displaystyle \frac{M}{n_0}}`$ (26)
$`\mathrm{\Gamma }_{m,0}={\displaystyle \frac{\mathrm{\Gamma }_m}{n_0^2}}`$ (27)
$`\gamma _0={\displaystyle \frac{\gamma }{n_0^2}}.`$ (28)
The effective Hamiltonian is:
$$_{\mathrm{e}ff}=\frac{3m^2}{8n_0^4}a_0(q)+\frac{E_0}{n_0^2}\gamma _0.$$
(29)
For a large microwave field amplitude or in the deep semiclassical limit, i.e. $`n_0\mathrm{}`$, we may employ the asymptotic expression, for large $`q`$, of the Mathieu parameter
$$a_0(q)2q+2\sqrt{q}.$$
(30)
This corresponds to the case where the pendulum is localized near its stable equilibrium position (whose energy is $`2q`$), its zero-point energy in the ground state being calculated in the harmonic approximation (hence the $`2\sqrt{q}`$ term).
In the opposite limit, i.e. for $`q1`$, another approximation exists $`a_0(q)q^2/2`$ . This corresponds basically to a very weak trapping pendulum potential, where the motion is basically the free motion only slightly perturbed (at second order in the energy) by the potential. This is the classical counterpart of the quantum perturbative approach developed above, the equivalent of the “no $`n`$-mixing” approximation. In the following, we restrict ourselves to this case as more appropriate for moderate $`n_0`$ and very small $`F_0`$. Hence, the final expression for the effective secular Hamiltonian we are going to deal with is
$$_{\mathrm{e}ff}=\frac{F_0^2}{3m^2}\mathrm{\Gamma }_{m,0}^2+\frac{E_0}{n_0^2}\gamma _0.$$
(31)
This Hamiltonian is a semiclassical counterpart of the quantum effective Hamiltonian, Eq. (14), namely first order in the static electric field and second order in the resonant microwave field.
To calculate the quasienergies semiclassically, one should quantize the secular motion determined by the Hamiltonian (31). It has been done in simpler situations (e.g. an hydrogen atom perturbed by a linearly polarized microwave field and a parallel static electric field ). Then, the secular motion is integrable and its quantization straightforward using the WKB quantization rule. The present secular problem has two degrees of freedom and turns out to be non-integrable except for some limiting cases. Therefore, a detailed analysis of the classical motion in the phase space of secular variables is necessary for possible comparison with quantal data.
## III PURE MICROWAVE PERTURBATION CASE
Let us consider first the perturbation of an hydrogen atom by an elliptically polarized microwave field in the absence of the static field. Our previous studies of resonant driving of the atom were restricted to the simplified two-dimensional model atom where the electronic motion is restricted to the polarization plane . Then, the classical secular motion is one-dimensional and application of the WKB quantization rule gives very accurate results for quasienergies of the system. In a realistic three-dimensional model, a similar procedure is no longer possible as the effective classical Hamiltonian for a general EP field is not integrable. The secular motion of the atom for weak microwaves, see Eq. (31), is determined by the Hamiltonian
$$_{\mathrm{e}ff}=\frac{F_0^2}{3m^2}\mathrm{\Gamma }_{m,0}^2(L_0,\mathrm{\Psi },M_0,\mathrm{\Phi };\alpha ).$$
(32)
The integrable motion is obtained in the limiting polarization cases, i.e. $`\alpha =0`$ or $`1`$, only. That is, for the LP field, the angular momentum projection on the polarization axis is a constant of motion. For the CP case, because of circular symmetry of the field, $`\mathrm{\Phi }`$ becomes a cyclic variable and the secular motion is also effectively one-dimensional. Clearly, for $`\alpha `$ close to these limiting values, the secular motion will be close to integrable. With this in mind, looking for chaotic dynamics, we take $`\alpha =0.4`$ in the following (we have verified that this value is representative for a typical EP behaviour).
Eq. (32) shows that the structure of the classical phase space of the secular motion depends only on the value of $`_{\mathrm{e}ff}/F_0^2`$ (beside the integer $`m`$ labeling the primary resonance). In other words, if the secular motion is non-integrable for some finite field amplitude, it remains non-integrable even for infinitesimally small amplitude. This clearly demonstrates the inapplicability of the Kolmogorov-Arnold-Moser theorem to the highly degenerate Coulomb problem. On the other hand, the time scale of precession of an electronic ellipse is affected by the strength of the perturbation; for very small $`F_0`$, it will be extremely slow, but the trajectories of the secular motion do not depend on $`F_0.`$
Consider the principal 1:1 resonance case, i.e. $`m=1`$. To focus on the phase space structure, we have plotted Poincaré surfaces of section (SOS) for a few values of $`_{\mathrm{e}ff}/F_0^2`$ in Fig. 1. For high values of the secular energy, the motion is generally regular. However, for an energy interval around $`_{\mathrm{e}ff}/F_0^2=0.3,`$ the mixed character of the motion is apparent, a quite large chaotic layer is clearly visible.
Switching to the 2:1 resonance case (right hand side of Fig. 1) we find, as previously, regular phase space structures for high energy and mixed character of the motion for lower energies. By visual inspection, the secular motion for $`2:1`$ resonance looks “more chaotic” with smaller regular islands embedded in a pronounced chaotic layer. Although we have searched quite extensively, we could not find values of $`\alpha `$ and the energy corresponding to fully chaotic motion. Always at best tiny regular islands have been found.
The mixed character of the secular motion should have consequences on the statistical properties of the quasienergy spectrum inside the $`n_0`$ manifold. As the system possesses anti-unitary symmetry invariance, i.e. is invariant under time-reversal combined with the $`yy`$ transformation, see Eq. (1), the statistical properties are expected to reflect intermediate behavior between the Poisson and GOE character. To calculate the quasienergy spectrum, we employ the quantum effective Hamiltonian Eq. (14). One should take care of discrete symmetries of the system. That is, the system is invariant under the $`zz`$ transformation as well as the parity combined with translation in time by $`\pi /\omega `$ transformations. Thus, the whole spectrum of the $`n_0`$-manifold splits into four independent smaller spectra which are unfolded independently in order to study level statistics.
The dynamics of the levels belonging to the $`n_0=20`$ manifold as a function of the polarization degree $`\alpha ,`$ are shown in Fig. 2, for the 1:1 and 2:1 resonance cases. In each panel, for clarity, there is only one of the four independent sub-spectra plotted. Qualitatively, the level dynamics reflects the character of classical motion, i.e. for high energies, one cannot see level repulsion and there are apparently level crossing (actually small avoided crossings). At lower energy, the level dynamics is more irregular with plenty of avoided crossings, a clear signature of classical chaos in the system, compare Fig. 3 which shows this region in more detail.
To make the comparison more quantitative, we have calculated the cumulative nearest neighbor spacing (NNS) distributions, taking levels in the energy intervals $`_{\mathrm{e}ff}/F_0^2[0.020.045]`$ and $`[0.000350.0018]`$ for the 1:1 and 2:1 resonance cases, respectively. The intervals have been chosen to correspond to chaotic behavior as much as possible. Fig. 4 presents the results for the 1:1 resonance for principal quantum number $`n_0`$ around 55 and for about twice bigger $`n_0`$, i.e. around 100. The similar results corresponding to the 2:1 resonance are plotted in Fig. 5. In all cases, one can observe that the numerical data are intermediate between the Poisson and GOE statistics. However, the behavior of the 1:1 case is closer to Poisson while the 2:1 one to GOE character, in agreement with the more classically chaotic behaviour in the latter case.
Quantitative measures can be obtained by fitting the data to some theoretical NNS distributions. Quite often, this is being done directly on the histograms for the NNS distributions. Such a procedure is, however, bin size dependent and should be avoided. We prefer to fit the integrated (cumulative) distributions which do not suffer from this ambiguity. There are several possible choices for the theoretical distributions. Berry-Robnik statistics corresponds to a superposition of independent Poisson and GOE spectra with relative weight $`q`$ corresponding to the relative volumes of regular and chaotic parts in the classical phase space. Others possibilities are the phenomenological Brody and Izrailev distributions. The explicit expressions for the distributions used can be found in Appendix B. Although the Berry-Robnik distribution relies on some reasonable theoretical grounds while the two other ones are purely phenomenological, it is commonly accepted that, for not very small effective $`\mathrm{}`$, the Brody distribution – or much less known Izrailev one – works better than the Berry-Robnik proposition. The latter works well only in the very deep semiclassical limit (very small effective $`\mathrm{}`$) where tunneling between regular and chaotic parts of the phase space is negligible. For lower lying states, the “regular” and “irregular” part of the spectrum interact via tunneling, leading to level repulsion between states in the two groups.
This behaviour is clearly visible in Figs. 4 and 5, (b) and (d). There, the Berry-Robnik distribution predicts much more small spacings than actually observed. On the other hand, except for small spacings we have found that Berry-Robnik statistics fits best in most of the cases, compare Fig. 4 and Fig. 5.
The obtained fitted values of the Berry-Robnik parameter and the parameters for the Brody and Izrailev distributions are given in Table I. While, for the latter cases, the parameters have little physical meaning, as mentioned above, $`q`$ in the Berry-Robnik distribution should measure the fraction of the chaotic phase space. Qualitatively, the value obtained agrees well with classical SOS plots.
We have also studied the spectral rigidity, i.e. the $`\mathrm{\Delta }_3`$ statistics , in order to obtain some information on the long range correlations in the spectra of our system. Spectral rigidity gives an independent information about the relative measure of the chaotic part of the classical phase space. For a superposition of independent Poisson and GOE spectra, one obtains
$$\mathrm{\Delta }_3=\mathrm{\Delta }_3^{\mathrm{𝑃𝑜𝑖𝑠𝑠𝑜𝑛}}((1q)L)+\mathrm{\Delta }_3^{\mathrm{𝐺𝑂𝐸}}(qL).$$
(33)
We have fitted our numerical results with this distribution. The results are plotted in Fig. 6 while values of the fitted parameter are put in Table I. It is well known that, at large $`L`$, the spectral rigidity deviates from any universal behaviour and saturates. This is a non-semiclassical effect which should take place for larger and larger $`L`$ as the effective $`\mathrm{}`$ goes to zero. To fit the parameter $`q`$, we have taken into account only data up to $`L=10`$ in order not to enter the saturation regions visible in Fig. 6.
Comparing the Poincaré SOS, NNS distributions and $`\mathrm{\Delta }_3`$ statistics of the calculated data, the qualitative agreement between the classical dynamics and the corresponding quantum statistical properties is apparent. The fitted values for the relative measure of the chaotic part of the phase space coming either from the NNS or the $`\mathrm{\Delta }_3`$ distributions agree perfectly and match well the visual aspect of the SOS.
The NNS distributions change a little when the principal quantum number is modified. The $`n_055`$ case reveals slightly stronger level repulsion (and, therefore, a “more chaotic” character) than the $`n_0100`$ case. This suggests that some tiny regular structures in the classical phase space are not resolved for $`n_0=55`$ but are for $`n_0=100.`$ The long range correlations are more dramatically sensitive to $`n_0`$. The saturation of $`\mathrm{\Delta }_3`$ starts at about twice larger distance in $`L`$ for the higher $`n_0`$ value, see Fig. 6. This is in agreement with the theory of Berry where the critical $`L`$ value scales as $`1/\mathrm{}`$ (the effective Planck constant in our problem is $`1/n_0`$). For large $`L,`$ the $`\mathrm{\Delta }_3`$ statistics saturates at almost twice bigger value for $`n_0100`$ than $`n_055`$ again in a qualitative agreement with theory. The latter predicts ($`\mathrm{\Delta }_3(\mathrm{})1/\mathrm{}`$) for a regular spectrum and ($`\mathrm{\Delta }_3(\mathrm{})\mathrm{ln}(1/\mathrm{})`$) for a strongly chaotic system. We deal with an intermediate system with mixed phase space, thus we expect the numerical $`\mathrm{\Delta }_3(\mathrm{})`$ value to lie in between the two limits. This is indeed the case.
Finally, let us briefly argue what happens for larger microwave field amplitudes. The secular motion, in the weak field limit, is determined by the effective Hamiltonian, Eq. (32), which in turn is a function of $`\mathrm{\Gamma }_{m,0}`$. Increasing the field amplitude one leaves the validity range of the Eq. (32) and enters the region where Eq. (30) is applicable. Then the orbital electronic motion becomes localized inside a resonance island. The secular motion, however, remains unchanged because the new effective Hamiltonian is again a function of $`\mathrm{\Gamma }_{m,0}`$ only. This means that the spectral statistical properties for higher field amplitude (of course not so big as to produce strong intermanifold mixing) will be the same as the ones presented here. For the high secular energy the motion will be also regular. Thus, one may expect that the nonspreading wavepackets predicted using the two-dimensional model will also exist in the real three-dimensional world.
## IV MICROWAVE PLUS STATIC ELECTRIC FIELD
In this section, we discuss the intramanifold behaviour of the hydrogen atom exposed to a resonant microwave field and a static electric field of arbitrary mutual orientation.
For small field amplitudes, the effective secular Hamiltonian is given by Eq. (31). The Hamiltonian is the sum of two terms – the first one proportional to $`F_0^2`$ (square of the scaled microwave field), the second one to $`E_0`$ (scaled static field). For arbitrary mutual orientations of the two fields, the two terms have incompatible symmetry properties and, when having comparable magnitudes, induce a globally chaotic behavior.
Eq. (31) has some well defined scaling properties with the field strengths $`F_0`$ and $`E_0`$. Let us define the reduced microwave strength
$$=F_0^2n_0^2/E_0=F^2n_0^6/E$$
(34)
and the reduced Hamiltonian
$$=_{\mathrm{e}ff}n_0^2/E_0=\frac{}{3m^2}\mathrm{\Gamma }_{m,0}^2+\gamma _0.$$
(35)
The classical phase space structure depends only on the value of $``$ and $``$ (beside the static field vector orientation and the polarization of the microwave field), but not on the detailed values of $`n_0`$, $`E_0`$, $`F_0`$ and the secular energy. Of course, weaker fields imply a slower secular motion, but this does not affect the structure of the phase space. In the quantum mechanical picture, the energy splitting of a degenerate hydrogenic manifold also depends on absolute values of the fields, but the structure of the levels does not.
The application of the static electric field allows us to break any anti-unitary symmetry of the system. It is only for $`E_x=0`$ (i.e. $`\phi =\pi /2`$) or $`E_y=0`$ (i.e. $`\phi =0`$) that the system has some anti-unitary invariance, under the combination of the time-reversal transformation with reflection with respect to the $`yz`$ or $`xz`$ plane, see Eq. (1).
In a previous letter , we have considered the system in the case of the 2:1 resonance driving, i.e. when the microwave frequency is twice the Kepler frequency. Let us consider here first the principal 1:1 resonance, i.e. $`m=1`$. A possible signature of the anti-unitary symmetry breaking would be to observe level repulsion with stronger than linear repulsion (i.e. with $`P(s)s^\beta `$ with $`\beta >1`$ for small $`s`$). Clearly, it is desirable to have a predominantly chaotic classical dynamics, as a transition from GOE to GUE statistics is expected. To this end, we have to find values of the fields parameters which maximize chaoticity of the system. After rather extensive searches, we have found that $`\alpha =0.4`$, $`=10`$ and $`\theta =\pi /4`$ are a good choice for that purpose. The remaining spherical angle $`\phi `$ determining the orientation of the static electric field vector, is used to control the breaking of the anti-unitary symmetry.
In Fig. 7, we show Poincaré SOS for a few values of $``$, for two cases: $`\phi =0`$ and $`\pi /2`$. One can see that, for $`\phi \pi /2`$, a predominantly chaotic structure space exists in phase space for a large range of secular energies. For $`\phi =\pi /2,`$ the generalized time reversal invariance holds, as mentioned above. Thus, to study the symmetry breaking, it is interesting to e.g. decrease (or increase) $`\phi `$ gradually, collecting quantum data for some $`\phi `$ values. Note that the addition of the static field tends to make the classical dynamics slightly more chaotic, compare Figs 1 and 7.
The data are then analyzed as in the previous Section – an example is shown in Fig. 8 for two different orientations of the electric field $`\phi =\{0.4\pi ,\pi /2\}`$. For each $`\phi `$, data have been collected for principal quantum number $`n_0`$ in the range 50-59. Then, we have chosen levels from the scaled energy interval $`(0,1.4)`$, unfolded each spectrum and calculated NNS distributions and spectral rigidities. The cumulative NNS distributions are plotted in Fig. 8. One can see that the NNS distribution corresponding to the anti-unitary invariant case, $`\phi =\pi /2`$, is close to, but does not reach completely the GOE behavior. Similarly, for the $`\phi =0.4\pi `$ case, the distribution is very close, but does not reach the GUE one. Nevertheless, the symmetry breaking is apparent and the numerical spectrum at $`\phi =0.4\pi `$ shows much more level repulsion than the GOE case, which is a clear-cut signature of the breaking of any anti-unitary symmetry.
To measure departures from the entirely ergodic behavior, we can also use the Berry-Robnik distribution. The Berry-Robnik model for fully broken antiunitary symmetry consists of the superposition of two independent Poisson and GUE spectra. The results of the fits are collected in Table II. We have not applied the Brody distribution to the broken anti-unitary invariance case, as the distribution is defined only for Brody parameter less than unity, and does not make any sense in the unitary case. The Izrailev distribution does not suffer this severe problem, and contain all limiting cases ($`\beta =0`$ for Poisson statistics, $`\beta =1`$ for GOE and $`\beta =2`$ for GUE). We have thus fitted our results with this distribution too.
In all cases, the Izrailev distribution works much better than the Berry-Robnik ansatz, presumably because chaotic motion occupies most of the phase space and regular regions are very tiny as seen from Fig. 7 and from the values of $`q`$ obtained.
The values of the fitted Berry-Robnik parameter are consistent with the character of the corresponding classical motion. An independent information about the relative measure of a chaotic part of phase space comes from the fit of the $`\mathrm{\Delta }_3`$ statistics, Fig. 9 \[for broken anti-unitary invariance case $`\mathrm{\Delta }_3^{\mathrm{𝐺𝑂𝐸}}`$ is substituted by $`\mathrm{\Delta }_3^{GUE}`$ in Eq. (33)\], which turns out to agree very nicely with the values of the Berry-Robnik parameter.
By gradually decreasing $`\phi ,`$ we may observe the partial symmetry breaking by studying the variation of fitted parameters with $`\phi `$. Such a transition has been analytically studied for gaussian random matrices where the two-point correlation function was analytically found for an appropriate ensemble which interpolated between the GOE and the GUE. A further link with the dynamics of fully chaotic systems with partially broken antiunitary symmetry was also established . These developments cannot be used here since the dynamics in our case is not fully chaotic (as seen from SOS plots and the non-integer level repulsion $`\beta `$ parameters for extreme cases of preserved and broken antiunitary symmetry, see Table II).
On the other hand, the Izrailev distribution is quite suitable since it should be a reasonable approximation both for a partial symmetry breaking and a mixed dynamics (it would be possible to construct an analog of Berry-Robnik distribution for such a case but it would be of little practical importance).
Fig. 10 summarizes the changes of the fitted Izrailev parameter $`\beta `$ (small $`s`$ repulsion) with $`\phi `$. For $`\phi =\pi /2,`$ it is minimal and equal to 0.85 (see also Table II) for $`n_0`$ around 55. With departure from the value $`\phi =\pi /2`$ where the anti-unitary symmetry exists, it rapidly increases (filled circles) up to the maximal value 1.73 for $`\phi =0.4\pi ,0.6\pi `$. For still lower (higher) values of $`\phi ,`$ the trend is reversed and $`\beta `$ starts to decrease. This, at first glance, is a surprising effect (since the symmetry should not yet start to restore). However, it is most probably due the fact that the classical motion becomes more regular as $`\phi `$ is far from $`\pi /2`$. Observe in Fig. 7 that the SOS around $`\phi =0`$ is much more regular than for $`\phi =\pi /2`$.
In order to test this hypothesis, we have used another $`n_0`$ value. Indeed, if the dip of $`\beta `$ near $`\phi =\pi /2`$ was due to classical reasons, it should not depend on $`\mathrm{}`$ ,i.e. $`n_0`$. Fig. 10 shows also the fitted $`\beta `$ parameters for $`n_0`$ around 100. Generally, the $`\beta `$ values obtained are slightly larger than for lower $`n_0`$. Clearly, $`\phi `$ starts to decrease again around the value $`\phi =0.4\pi `$. On the other hand, the dip of $`\beta `$ when $`\phi `$ goes to $`\pi /2`$ (where an anti-unitary symmetry exists) is faster than for lower $`n_0`$. This is in a full qualitative agreement with RMT : for the GOE $``$ GUE transition the parameter controlling the transition is proportional to $`N^{1/2}`$ where $`N`$ is the matrix dimension. The size of our matrices scales as $`n_0^2`$, so the parameter controlling the breaking (i.e. a deviation from $`\pi /2`$ value) should scale like $`n_0^1\mathrm{}_{\mathrm{e}ff}`$. Such a behaviour is roughly observed in Fig.10.
Fig. 11 shows the “maximal” repulsion case obtained, i.e. data for $`n_0`$ around 100 and $`\phi =0.4\pi `$. The numerical data are presented in the form of the histogram of spacings and compared with the Izrailev distribution (the fit has been performed, as usual, for the cumulative distribution; the resulting $`\beta =1.83`$ value has been used to plot the Izrailev distribution). The dash-dotted and dashed curves correspond to Wigner GOE and GUE distributions, respectively. The fact that we observe level repulsion much stronger than the GOE behaviour is a signature of anti-unitary symmetry breaking.
As mentioned before, our first results on the manifestations of symmetry breaking in our system have been obtained for the microwave frequency being twice the Kepler frequency, i.e. for the 2:1 resonant driving . This choice was motivated by SOS plots in the absence of the electric field - see Fig. 1 \- showing smaller regular islands for higher microwave frequency. However, as we have seen above, the presence of the electric field makes the secular motion in the principal resonance island chaotic enough and in fact we get stronger repulsion for the $`1:1`$ resonance than for the $`2:1`$ situation reported before . For completeness, we show in Table III the fitted parameters obtained from the numerical results reported in using either the spacing distribution or the spectral rigidity $`\mathrm{\Delta }_3.`$ As can be seen, the conclusions we obtain from these results completely confirm the analysis of the 1:1 resonance.
As far as we know, the studied system constitutes the first experimentally realizable example of a quantum system with broken anti-unitary symmetry. We have considered small, but finite, field amplitudes to stay well within the applicability range of the perturbation theory. Nevertheless, this is experimentally feasible: for $`n_055`$ and $`F_05\times 10^4,`$ i.e. about $`0.3`$ V/cm, the mean level spacing is of the order of MHz.
For stronger fields, our classical studies also suggest a similar behavior. It could even be that the breaking of the secular approximation makes the system more chaotic and that the statistical properties will be closer to GUE. However, we are not able to show quantum numerical results as they require full quantum numerical treatment which is difficult with the present computer resources and must be left for a future work.
## V CONCLUSIONS
We have considered an hydrogen atom perturbed by a resonant elliptically polarized microwave field with or without an additional static electric field of different orientation, in the limit of small field amplitudes. Classically, such fields may produce chaotic dynamics in the secular motion of the electronic elliptical trajectory. In quantum mechanical language, states coming from a given hydrogenic manifold may be mixed significantly only with each other. Such a situation has been interpreted as a signature of an intramanifold chaos.
For the pure microwave problem, we have studied two different resonant driving cases, i.e. the 1:1 and 2:1 resonances between the microwave field and the unperturbed Kepler motion. Quantizing the fast orbital electronic motion, one can derive an effective Hamiltonian describing the slow secular precession of the electronic elliptical trajectory. For a generic elliptical polarization, the effective Hamiltonian has two degrees of freedom and turns out to be non-integrable. By means of Poincaré surfaces of sections, we have found a range of the secular energy where the phase space reveals mixed character with a significant amount of chaotic layer. Switching to a quantum perturbation calculation, we have shown that the statistical properties of the corresponding quasienergy levels reveal an intermediate behavior between Poisson and GOE character. For the 2:1 resonance case, the classical phase space is significantly more irregular than for the 1:1 case and, consequently, the spectral properties are closer to the GOE behavior.
The application of an additional static electric field to the system has allowed us to enhance chaos in the secular motion. Moreover, the static electric field, for a generic orientation, breaks any anti-unitary symmetry of the system which has a dramatic effect on statistical properties of quasienergy levels. This is the first, to our knowledge, experimentally realizable quantum system, with corresponding chaotic classical behavior, which exhibits breaking of any generalized time-reversal symmetry.
We have studied the principal 1:1 resonance for two slightly different static field orientations: the first one corresponds to the anti-unitary invariance case, the other one to breaking such a symmetry. The classical phase space structures, in both cases, are similar with predominately chaotic behavior. However, the statistical properties of the quantum spectrum change from a near-GOE to a near-GUE behavior when one switches from the preserved to broken anti-unitary symmetry case. In the intermediate situations, we could observe a partial symmetry breaking effect due to the finite size of matrices involved in the problem.
We have studied the limit of small field amplitudes as it allows us to employ quantum perturbation theory. For higher amplitudes, the classical effective Hamiltonian is known, but finding the quasienergy spectrum requires full quantum numerical calculations. For a pure microwave perturbation, however, we may predict that statistical properties of the quantum spectrum, for stronger field, should be the same as for weak field limit. This is because of the specific form of the effective Hamiltonian, Eq. (21), which depends on dynamical variables only through the $`\mathrm{\Gamma }_m`$.
## VI Acknowledgements
We are grateful to Felix Izrailev for the permission to use his code for his NNS distribution. We thank the referee for bringing references to our attention. Support of KBN under project 2P302B-00915 (KS and JZ) is acknowledged. KS acknowledges support by the Alexander von Humboldt Foundation. The additional support of the bilateral Polonium and PICS programs is appreciated. Laboratoire Kastler Brossel de l’Université Pierre et Marie Curie et de l’Ecole Normale Supérieure is UMR 8552 du CNRS.
## VII APPENDIX A
The amplitude and the phase of the atomic dipole at the microwave frequency, $`\mathrm{\Gamma }_m(L,\mathrm{\Psi },M,\mathrm{\Phi };\alpha )`$ and $`\delta (L,\mathrm{\Psi },M,\mathrm{\Phi };\alpha )`$ appearing in Eq. (16) may be expressed as
$`\mathrm{\Gamma }_m`$ $`=`$ $`\{\left({\displaystyle \frac{1+\alpha }{2}}\right)^2[V_m^2+U_m^2]+\left({\displaystyle \frac{1\alpha }{2}}\right)^2[V_m^2+U_m^2]`$ (37)
$`+{\displaystyle \frac{1\alpha ^2}{2}}[(V_mV_mU_mU_m)\mathrm{cos}2\mathrm{\Phi }(V_mU_m+U_mV_m)\mathrm{sin}2\mathrm{\Phi }]\}^{1/2},`$
and
$$\mathrm{tan}\delta =\frac{(1\alpha )(V_m\mathrm{sin}\mathrm{\Phi }+U_m\mathrm{cos}\mathrm{\Phi })(1+\alpha )V_m\mathrm{sin}\mathrm{\Phi }+U_m\mathrm{cos}\mathrm{\Phi })}{(1\alpha )(V_m\mathrm{cos}\mathrm{\Phi }U_m\mathrm{sin}\mathrm{\Phi })+(1+\alpha )V_m\mathrm{cos}\mathrm{\Phi }U_m\mathrm{sin}\mathrm{\Phi })},$$
(38)
where $`V_m`$ and $`U_m`$ are Fourier expansion terms of the original hamiltonian (1) in action-angle variables. Explicitly, they read
$`V_0(J,L,\mathrm{\Psi })`$ $`=`$ $`{\displaystyle \frac{3e}{2}}J^2\mathrm{cos}\mathrm{\Psi },`$ (39)
$`U_0(J,L,M,\mathrm{\Psi })`$ $`=`$ $`{\displaystyle \frac{3eM}{2L}}J^2\mathrm{sin}\mathrm{\Psi },`$ (40)
and for $`m0`$
$`V_m(J,L,M,\mathrm{\Psi })`$ $`=`$ $`{\displaystyle \frac{J^2}{m}}[𝒥_m^{^{}}(me)+{\displaystyle \frac{M\sqrt{1e^2}}{Le}}𝒥_m(me)]\mathrm{cos}\mathrm{\Psi },`$ (41)
$`U_m(J,L,M,\mathrm{\Psi })`$ $`=`$ $`{\displaystyle \frac{J^2}{m}}[{\displaystyle \frac{M}{L}}𝒥_m^{^{}}(me)+{\displaystyle \frac{\sqrt{1e^2}}{e}}𝒥_m(me)]\mathrm{sin}\mathrm{\Psi }.`$ (42)
In the above formulae $`e=\sqrt{1L^2/J^2}`$ is an eccentricity of the electronic ellipse while $`𝒥_m(x)`$ and $`𝒥_m^{^{}}(x)`$ denote the ordinary Bessel function and its derivative, respectively.
## VIII APPENDIX B
We give here explicit expressions for various level spacing distributions which have been used in the present paper.
The Poisson distribution, corresponding to a system with classically integrable dynamics , reads
$$P(s)=\mathrm{exp}\left(s\right).$$
(43)
For an ergodic classical behavior, the quantum spectrum of a generic system is conjectured to have a nearest neighbor spacing (NNS) distribution (for the unfolded spectrum) similar to that of the random matrices of the same universality class. The resulting NNS distributions are quite complicated (see e.g. ). However, a good approximation is given by the so called Wigner surmise, obtained for matrices of rank 2. These are:
$$P(s)=\frac{\pi }{2}s\mathrm{exp}\left(\frac{\pi }{4}s^2\right)$$
(44)
for an anti-unitary invariant (GOE) system and
$$P(s)=\frac{32}{\pi ^2}s^2\mathrm{exp}\left(\frac{4}{\pi }s^2\right)$$
(45)
for broken anti-unitary invariance (GUE).
The phenomenological Brody distribution which interpolates between Poisson and GOE distributions reads
$$P(s)=C(\beta +1)s^\beta \mathrm{exp}\left(Cs^{\beta +1}\right)$$
(46)
with
$$C=\left[\mathrm{\Gamma }\left(\frac{\beta +2}{\beta +1}\right)\right]^{\beta +1},$$
(47)
$`\beta =0`$ (resp. $`\beta =1`$) corresponds to the extreme case of the Poisson (resp. GOE) statistics.
Another attempt towards a simple distribution interpolating between different ensembles, is due to Izrailev and reads
$$P(s)=As^\beta (1+sB\beta )^{f(\beta )}\mathrm{exp}\left[\frac{\pi ^2\beta s^2}{16}\frac{\pi }{4}(2\beta )s\right]$$
(48)
where
$$f(\beta )=\frac{2^\beta (1\beta /2)}{\beta }0.16874,$$
(49)
and $`A`$, $`B`$ are constants that ensure
$$P(s)𝑑s=1,$$
(50)
and
$$sP(s)𝑑s=1.$$
(51)
It is claimed to work reasonably well for all possible intermediate situations.
While there exist several other propositions in the literature, we list only the so called Berry–Robnik statistics . It may be derived assuming an independent superposition of Poisson spectrum and spectra corresponding to random matrix predictions. If one deals with Poisson and only one GOE spectrum, the Berry–Robnik distribution reads
$$P(s)=\left[2q(1q)+\frac{\pi }{2}q^3s\right]\mathrm{exp}\left[(q1)s\frac{\pi }{4}q^2s^2\right]+(1q)^2\mathrm{exp}[(q1)s]\text{erfc}\left(\frac{\sqrt{\pi }}{2}qs\right)$$
(52)
where $`0q1`$ is the relative weight of the GOE spectrum. Classically, $`q`$ corresponds to the relative volume of the chaotic part of phase space. The similar Berry–Robnik distribution for a superposition of a Poisson spectrum and one GUE spectrum is
$`P(s)=\{[{\displaystyle \frac{32}{\pi ^2}}q^4s^2+{\displaystyle \frac{8}{\pi }}q^2(1q)s+(1q)^2]\mathrm{exp}({\displaystyle \frac{4}{\pi }}q^2s^2)`$ (53)
$`+[2q(1q)(1q)^2qs]\text{erfc}\left({\displaystyle \frac{2}{\sqrt{\pi }}}qs\right)\left\}\mathrm{exp}\right[(q1)s].`$ (54)
Let us present also the expression for the $`U(W)`$ function which has been employed in a fine-scale representation of the deviation of the numerical level spacing distribution from the best fitting theoretical distribution. The following function
$$U(W)=\mathrm{arccos}\sqrt{1W},$$
(55)
where $`W`$ is the value of the cumulative level spacing distribution $`P(s)𝑑s`$, ensures that, over the whole range of $`W`$, i.e. from 0 to 1, the standard deviation of numerical data is uniform and equal to $`\delta U=1/(\pi \sqrt{N})`$, where $`N`$ is the total number of spacings.
For completeness, let us finally define the spectral rigidity $`\mathrm{\Delta }_3(L)`$ as an average over the spectral range used in analysis (i.e. over $`x_0`$) of
$$\mathrm{\Delta }_3(x_0,L)=L^1\mathrm{m}in_{A,B}_{x_0}^{x_0+L}𝑑x(N(x)AxB)^2,$$
(56)
where $`N(x)`$ is the integrated level density (a staircase function). |
warning/0002/math0002050.html | ar5iv | text | # Minimal submanifolds of Kähler-Einstein manifolds with equal Kähler angles
## 1 Introduction
Let $`(N,J,g)`$ be a Kähler manifold of complex dimension $`2n`$ and $`F:MN`$ an immersed submanifold of real dimension $`2n`$. We denote by $`\omega `$ the Kähler form of $`N`$, $`\omega (X,Y)=g(JX,Y)`$. On $`M`$ we take the induced metric $`g_M=F^{}g`$. $`N`$ is Kähler-Einstein if its Ricci tensor is a multiple of the metric, $`Ricci^N=Rg`$. At each point $`pM`$, we identify $`F^{}\omega `$ with a skew-symmetric operator of $`T_pM`$ by using the musical isomorphism with respect to $`g_M`$, namely $`g_M(F^{}\omega (X),Y)=F^{}\omega (X,Y)`$. We take its polar decomposition
$$F^{}\omega =\stackrel{~}{g}J_\omega $$
(1.1)
where $`J_\omega :T_pMT_pM`$ is a ( in fact unique) partial isometry with the same kernel $`𝒦_\omega `$ as of $`F^{}w`$, and where $`\stackrel{~}{g}`$ is the positive semidefinite operator $`\stackrel{~}{g}=|F^{}\omega |=\sqrt{(F^{}\omega )^2}`$. It turns out that $`J_\omega :𝒦_\omega ^{}𝒦_\omega ^{}`$ defines a complex structure on $`𝒦_\omega ^{}`$, the orthogonal compliment of $`𝒦_\omega `$ in $`T_pM`$. Moreover, it is $`g_M`$-orthogonal. If we denote by $`\mathrm{\Omega }_{2k}^0`$ the largest open set of $`M`$ where $`F^{}\omega `$ has constant rank $`2k`$, $`0kn`$, then $`𝒦_\omega ^{}`$ is a smooth sub-vector bundle of $`TM`$ on $`\mathrm{\Omega }_{2k}^0`$. Moreover, $`\stackrel{~}{g}`$ and $`J_\omega `$ are both smooth on these open sets. The tensor $`\stackrel{~}{g}`$ is continuous on all $`M`$ and locally Lipschitz, for the map $`P|P|`$ is Lipschitz in the space of normal operators. Let $`\{X_\alpha ,Y_\alpha \}_{1\alpha n}`$ be a $`g_M`$-orthonormal basis of $`T_pM`$, that diagonalizes $`F^{}\omega `$ at $`p`$, that is
$$F^{}\omega =\underset{1\alpha n}{}\left[\begin{array}{cc}0& \mathrm{cos}\theta _\alpha \\ \mathrm{cos}\theta _\alpha & 0\end{array}\right],$$
(1.2)
where $`\mathrm{cos}\theta _1\mathrm{cos}\theta _2\mathrm{}\mathrm{cos}\theta _n0`$. The angles $`\{\theta _\alpha \}_{1\alpha n}`$ are the Kähler angles of $`F`$ at $`p`$. Thus, $`\alpha `$, $`F^{}\omega (X_\alpha )=\mathrm{cos}\theta _\alpha Y_\alpha `$, $`F^{}\omega (Y_\alpha )=\mathrm{cos}\theta _\alpha X_\alpha `$ and if $`k1`$, where $`2k`$ is the rank of $`F^{}\omega `$ at $`p`$, $`J_\omega X_\alpha =Y_\alpha `$ $`\alpha k`$. The Weyl’s perturbation theorem applied to the eigenvalues of the symmetric operator $`|F^{}\omega |`$ shows that, ordering the $`\mathrm{cos}\theta _\alpha `$ in the above way, the map $`pcos\theta _\alpha (p)`$ is locally Lipschitz on $`M`$, for each $`\alpha `$. A complex direction of $`F`$ is a real two-plane $`P`$ of $`T_pM`$ such that $`dF(P)`$ is a complex line of $`T_{F(p)}N`$, i.e., $`JdF(P)dF(P)`$. Similarly, $`P`$ is said to be a Lagrangian direction of $`F`$ if $`\omega `$ vanishes on $`dF(P)`$, that is, $`JdF(P)dF(P)`$. The immersion $`F`$ has no complex directions iff $`\mathrm{cos}\theta _\alpha <1`$ $`\alpha `$. $`M`$ is a complex submanifold iff $`cos\theta _\alpha =1`$ $`\alpha `$, and is a Lagrangian submanifold iff $`cos\theta _\alpha =0`$ $`\alpha `$. We say that $`F`$ has equal Kähler angles if $`\theta _\alpha =\theta `$ $`\alpha `$. Complex and Lagrangian submanifolds are examples of such case. If $`F`$ is a complex submanifold, then $`J_\omega `$ is the complex structure induced by $`J`$ of $`N`$. The Kähler angles are some functions that at each point $`p`$ of $`M`$ measure the deviation of the tangent plane $`T_pM`$ of $`M`$ from a complex or a Lagrangian subspace of $`T_{F(p)}N`$. This concept was introduced by Chern and Wolfson \[Ch-W\] for surfaces, namely $`F^{}\omega =\mathrm{cos}\theta Vol_M`$. This $`\mathrm{cos}\theta `$ may have negative values and is smooth on all $`M`$. In our definition, for $`n=1`$, we demanded $`\mathrm{cos}\theta 0`$, that is, it is the modulus of the $`\mathrm{cos}\theta `$ given for surfaces. This may make our $`\mathrm{cos}\theta `$ do not be smooth. We have chosen this definition, because in higher dimensions we do not have a preferential orientation assigned to the real planes $`span\{X_\alpha ,Y_\alpha \}`$.
Our main aim is to find conditions for a minimal submanifold $`F`$ to be Lagrangian or complex, or $`M`$ to be a Kähler manifold with respect to $`J_\omega `$. The first result in this direction is due to Wolfson, for the case $`n=1`$:
###### Theorem 1.1
\[W\] If $`M`$ is a real compact surface and $`N`$ is a complex Kähler-Einstein surface with $`R<0`$, anf if $`F`$ is minimal with no complex points, then $`F`$ is Lagrangian.
Some results of \[S-V\] are a generalization of the above theorem to higher dimensions. In this paper we study the case of equal Kähler angles. Let us denote by $`_XdF(Y)=_{}dF(X,Y)`$ the second fundamental form of $`F`$. It is a symmetric tensor and takes values in the normal bundle $`NM=(dF(TM))^{}`$. $`F`$ is minimal iff $`trace_{g_M}_{}dF=0`$. Let $`()^{}`$ denote the orthogonal projection of $`F^1TN`$ onto the normal bundle. If $`F`$ is an immersion with no complex directions at $`p`$ and $`\{X_\alpha ,Y_\alpha \}`$ diagonalizes $`F^{}\omega `$ at $`p`$, then $`\{dF(Z_\alpha )`$, $`dF(Z_{\overline{\alpha }})`$, $`(JdF(Z_\alpha ))^{}`$, $`(JdF(Z_{\overline{\alpha }}))^{}\}`$ constitutes a complex basis of $`T_{F(p)}^cN`$, where
$$Z_\alpha =\frac{X_\alpha iY_\alpha }{2}=\mathrm{`}\mathrm{`}\alpha \mathrm{"},Z_{\overline{\alpha }}=\overline{Z_\alpha }=\frac{X_\alpha +iY_\alpha }{2}=\mathrm{`}\mathrm{`}\overline{\alpha }\mathrm{"}$$
(1.3)
are complex vectors of the complexified tangent space of $`M`$ at $`p`$. We extend to the complexified vector bundles the Riemannian tensor metric $`g_M`$ (sometimes denoted by $`,`$), the curvature tensors of $`M`$ and $`N`$, and any other tensors that will occur, always by $`lC`$-multilinearity. On $`M`$ the Ricci tensor of $`N`$ can be described by the following expression (\[S-V\]): for $`U,VT_{F(p)}N`$,
$$Ricci^N(U,V)=\underset{1\mu n}{}\frac{4}{\mathrm{sin}^2\theta _\mu }R^N(U,JV,dF(\mu ),(JdF(\overline{\mu }))^{}),$$
(1.4)
where $`R^N`$ denotes the Riemannian curvature tensor of $`N`$. An application of Codazzi equation to the above expression proves that, if $`N`$ is Kähler-Einstein with $`R0`$, Theorem 1.1 can be generalized to any dimension for totally geodesic maps (\[S-V\]).
We can also obtain the same conclusion to “broadly-pluriminimal” immersions for $`n=2`$, and $`N`$ Kähler-Einstein with negative Ricci tensor (\[S-V\]). A minimal immersion $`F`$ is said to be broadly-pluriminimal, if, for each $`p\mathrm{\Omega }_{2k}^0`$, with $`k1`$, $`F`$ is pluriharmonic with respect to any $`g_M`$-orthogonal complex structure $`\stackrel{~}{J}=J_\omega J^{}`$ on $`T_pM`$ where $`J^{}`$ is any $`g_M`$-orthogonal complex structure of $`𝒦_\omega `$ at $`p`$, that is, $`(_{}dF)^{(1,1)}=0`$. The (1,1)-part of $`_{}dF`$ is just given by $`(_{}dF)^{(1,1)}(X,Y)=\frac{1}{2}(_{}dF(X,Y)+_{}dF(\stackrel{~}{J}X,\stackrel{~}{J}Y))X,YT_pM.`$ If $`𝒦_\omega =0`$, this means that $`F`$ is pluriharmonic with respect to the almost complex structure $`J_\omega `$ (see for example \[O-V\]). In this case, we say that $`F`$ is pluriminimal in the usual sense, or simply pluriminimal. Pluriharmonic immersions are obviously minimal. If $`F`$ has equal Kähler angles, then only $`\mathrm{\Omega }_{2n}^0`$ is considered, where $`𝒦_\omega =0`$ and $`\stackrel{~}{J}=J_\omega `$. Products of minimal real surfaces of Kähler surfaces, totally geodesic submanifolds, minimal Lagrangian submanifolds, and complex submanifolds are examples of broadly-pluriminimal submanifolds. We will see in sections 2 and 3 that the concept of broadly-pluriminimality, for immersions without complex directions and with equal Kähler angles, may have a geometric interpretation in terms of the torsion of a new Riemannian connection on $`TM`$, described through an isomorphism $`\mathrm{\Phi }`$ from the tangent bundle of $`M`$ into the normal bundle. Pluriminimal immersions with equal Kähler angles immersed into Kähler-Einstein manifolds, that are not complex submanifolds, have constant Kähler angle, and only exist on Ricci-flat manifolds. In this case, $`\mathrm{\Phi }`$ defines a parallel homothetic isomorphism between $`TM`$ and $`NM`$.
For a minimal immersion $`F`$ with no complex directions we consider the locally Lipschitz map, symmetric on the Kähler angles,
$$\kappa =\underset{1\alpha n}{}log\left(\frac{1+\mathrm{cos}\theta _\alpha }{1\mathrm{cos}\theta _\alpha }\right).$$
(1.5)
This map is smooth on each $`\mathrm{\Omega }_{2k}^0`$, non-negative, and vanishes at Lagrangian points. It is an increasing map on each $`\mathrm{cos}\theta _\alpha `$. In \[S-V\] we have given an expression for $`\mathrm{}\kappa `$ at a point $`p_0\mathrm{\Omega }_{2k}^0`$, which we prove in the appendix of this paper, namely,
$`\mathrm{}\kappa `$ $`=`$ $`4i{\displaystyle \underset{\beta }{}}Ricci^N(JdF(\beta ),dF(\overline{\beta }))`$
$`+{\displaystyle \underset{\beta ,\mu }{}}{\displaystyle \frac{32}{\mathrm{sin}^2\theta _\mu }}Im(R^N(dF(\beta ),dF(\mu ),dF(\overline{\beta }),JdF(\overline{\mu })+i\mathrm{cos}\theta _\mu dF(\overline{\mu })))`$
$`{\displaystyle \underset{\beta ,\mu ,\rho }{}}{\displaystyle \frac{64(\mathrm{cos}\theta _\mu +\mathrm{cos}\theta _\rho )}{\mathrm{sin}^2\theta _\mu \mathrm{sin}^2\theta _\rho }}Re\left(g(_\beta dF(\mu ),JdF(\overline{\rho }))g(_{\overline{\beta }}dF(\rho ),JdF(\overline{\mu }))\right)`$
$`+{\displaystyle \underset{\beta ,\mu ,\rho }{}}{\displaystyle \frac{32(\mathrm{cos}\theta _\rho \mathrm{cos}\theta _\mu )}{\mathrm{sin}^2\theta _\mu \mathrm{sin}^2\theta _\rho }}(|g(_\beta dF(\mu ),JdF(\rho ))|^2+|g(_{\overline{\beta }}dF(\mu ),JdF(\rho ))|^2)`$
$`+{\displaystyle \underset{\beta ,\mu ,\rho }{}}{\displaystyle \frac{32(\mathrm{cos}\theta _\mu +\mathrm{cos}\theta _\rho )}{\mathrm{sin}^2\theta _\mu }}\left(|_\beta \mu ,\rho |^2+|_{\overline{\beta }}\mu ,\rho |^2\right),`$
where $`\{X_\alpha ,Y_\alpha \}_{1\alpha n}`$ is a $`g_M`$-orthonormal local frame of $`M`$, with $`Y_\alpha =J_\omega X_\alpha `$ for $`\alpha k`$, $`\{X_\alpha ,Y_\alpha \}_{\alpha k+1}`$ any $`g_M`$-orthonormal frame of $`𝒦_\omega `$, and which at $`p_0`$ diagonalizes $`F^{}\omega `$. For $`F`$ pluriminimal on $`\mathrm{\Omega }_{2n}^0`$ and $`N`$ Kähler-Einstein , we can get the following very simple final expression on $`\mathrm{\Omega }_{2n}^0`$ (\[S-V\])
$$\mathrm{}\kappa =2R\left(\underset{1\beta n}{}\mathrm{cos}\theta _\beta \right).$$
(1.7)
If $`F`$ has equal Kähler angles, then the expression of $`\mathrm{}\kappa `$ given in (1.6) can also be substantially simplified. Minimal surfaces with constant curvature and constant Kähler angle in complex space forms have been classified in \[O\]. Conditions on the curvature of $`M`$, $`N`$, and/or constant equal Kähler angles lead to some conclusions in our case as well, as we show in the theorems below. Henceforth, we assume $`N`$ is Kähler-Einstein. The expression for $`\mathrm{}\kappa `$, where the Ricci tensor of $`N`$ appears, and the Weitzenböck formula for $`F^{}\omega `$, leading to an integral equation involving the scalar curvature $`R`$, some trigonometric functions of the common Kähler angle, and the gradient of its cosine (Proposition 4.2), are our tools to obtain the results of this paper. In section 4 we prove our main results, namely:
###### Theorem 1.2
Let $`F`$ be a minimal immersion of a compact oriented manifold $`M`$, into a Kähler-Einstein manifold $`N`$, with equal Kähler angles.
(i) If $`n=2`$ and $`R0`$, then $`F`$ is either a complex or a Lagrangian submanifold.
(ii) If $`n3`$, $`R<0`$, and $`F`$ has no complex points, then $`F`$ is Lagrangian.
(iii) If $`n3`$, $`R=0`$, and $`F`$ has no complex points, then the common Kähler angle must be constant.
The conclusions in $`(i)`$ and $`(ii)`$ give a generalization of Theorem 1.1 to higher dimensions and equal Kähler angles. The case $`n=2`$ is the most special, because, in this dimension, immersions with equal Kähler angles have harmonic $`F^{}\omega `$, as we will see in section 3. The case $`n=3`$ also has special properties. If the angle is constant we may allow $`R>0`$:
###### Theorem 1.3
Let $`F`$ be minimal with constant equal Kähler angles, $`M`$ compact, orientable, and $`R0`$. Then, $`F`$ is either a complex or a Lagrangian submanifold.
###### Theorem 1.4
Let $`F`$ be minimal with equal Kähler angles, and M compact, orientable, with non-negative isotropic scalar curvature. If $`n=2`$ or $`3`$, then one of the following cases holds:
$`(i)`$ $`M`$ is a complex submanifold of $`N`$.
$`(ii)`$ $`M`$ is a Lagrangian submanifold of $`N`$.
$`(iii)`$ $`R=0`$ and $`\mathrm{cos}\theta =constant0,1`$, $`J_\omega `$ is a complex integrable structure, with $`(M,J_\omega ,g_M)`$ a Kähler manifold.
For any $`n1`$, any $`R`$, and constant equal Kähler angle, $`(i),(ii)`$ or $`(iii)`$ hold as well.
This theorem can be applied, for instance, to flat minimal tori on Calabi-Yau manifolds, or to spheres or products of $`S^2`$ with $`S^2`$ or with flat tori minimaly immersed into Kähler-Einstein manifolds with positive scalar curvature.
## 2 The morphism $`\mathrm{\Phi }`$
We consider the following morphism of vector bundles
$$\begin{array}{cccc}\mathrm{\Phi }:& TM& & NM\\ & X& & (JdF(X))^{}\end{array}$$
We easily verify that
$$\mathrm{\Phi }(X)=JdF(X)dF(F^{}\omega (X)).$$
(2.1)
Both $`TM`$ and $`NM`$ are real vector bundles of the same dimension $`2n`$. $`F`$ has no complex directions iff $`\mathrm{\Phi }`$ is an isomorphism. In fact $`\mathrm{\Phi }(X)=0`$, iff $`JdF(X)=dF(Y)`$ for some $`Y`$, i.e., $`span\{X,Y=\mathrm{`}\mathrm{`}JX\mathrm{"}\}`$ is a complex direction of $`F`$. Assume there are no complex directions. Then,
$$\widehat{g}(X,Y)=g_M(X,Y)g_M(F^{}\omega (X),F^{}\omega (Y))$$
(2.2)
defines a Riemannian metric on $`M`$. With this metric, $`\mathrm{\Phi }:(TM,\widehat{g})(NM,g)`$ is an isomorphism of Riemannian vector bundles. Let us denote by $`_{}`$, $`\widehat{}_{}`$, $`_{}^{}`$, and $`_{}^{}`$, respectively, the Levi-Civita connection of $`(M,g_M)`$, the Levi-Civita connection of $`(M,\widehat{g})`$, the usual connection of $`NM`$ induced by the Levi-Civita connection of $`N`$, and the connection on $`TM`$ that makes the isomorphism $`\mathrm{\Phi }`$ parallel, namely $`_{}^{}=\mathrm{\Phi }^1_{}^{}`$. We will also denote by $`_{}`$ the Levi-Civita connection of $`N`$ and the induced connection on $`F^1TN`$, as well. Thus, if $`U`$ is a smooth section of $`NMF^1TN`$, and $`X,Y`$ are smooth vector fields on $`M`$, we have
$$_X^{}U=(_XU)^{}\mathrm{\Phi }(_X^{}Y)=_X^{}(\mathrm{\Phi }(Y)).$$
The connections $`_{}`$ and $`\widehat{}_{}`$ have no torsion, because they are Levi-Civita, but $`_{}^{}`$ may have non-zero torsion $`T^{}`$. Since both $`\widehat{}_{}`$ and $`_{}^{}`$ are Riemannian connections of $`TM`$ for the same Riemannian metric $`\widehat{g}`$, then $`T^{}=0`$ iff $`\widehat{}_{}=_{}^{}`$ iff $`\mathrm{\Phi }`$ is parallel. Note that, if $`F`$ is Lagrangian, then $`\mathrm{\Phi }(X)=JdF(X)NM`$, $`J(NM)=dF(TM)`$, and $`\widehat{g}=g_M`$, $`\widehat{}_{}=_{}`$. Therefore, $`_X\mathrm{\Phi }(Y)=\left(J_XdF(Y)\right)^{}=0`$, that is, $`\mathrm{\Phi }`$ is parallel, and so $`_{}^{}=_{}`$, as well. In the next section (Corollary 3.2), we will see a converse of this. We extend $`\mathrm{\Phi }:TM^cNM^c`$ to the complexified spaces by $`lC`$-linearity.
###### Lemma 2.1
If $`\{X_\alpha ,Y_\alpha \}`$ is a diagonalizing $`g_M`$-orthonormal basis of $`F^{}\omega `$ at $`p`$, then at $`p`$, and for each $`\alpha ,\beta `$
$`\mathrm{\Phi }(T^{}(Z_\alpha ,Z_{\overline{\beta }}))`$ $`=`$ $`i(\mathrm{cos}\theta _\alpha +\mathrm{cos}\theta _\beta )_{Z_\alpha }dF(Z_{\overline{\beta }})`$
$`\mathrm{\Phi }(T^{}(Z_\alpha ,Z_\beta ))`$ $`=`$ $`i(\mathrm{cos}\theta _\alpha \mathrm{cos}\theta _\beta )_{Z_\alpha }dF(Z_\beta ).`$
Proof.
$`\mathrm{\Phi }(_X^{}Y)`$ $`=`$ $`_X^{}(\mathrm{\Phi }(Y))=(_X(\mathrm{\Phi }(Y)))^{}=(_X(JdF(Y)dF(F^{}\omega (Y))))^{}`$
$`=`$ $`\left(J_XdF(Y)+JdF(_XY)_XdF(F^{}\omega (Y))\right)^{}.`$
Therefore, using the symmetry of the $`_{}dF`$ and the fact that $`_{}`$ is torsionless,
$$\mathrm{\Phi }(T^{}(X,Y))=\mathrm{\Phi }(_X^{}Y_Y^{}X[X,Y])=_XdF(F^{}\omega (Y))+_YdF(F^{}\omega (X)).$$
(2.3)
The lemma follows now immediately. $`\mathbf{}`$
For each $`UNM_p`$, let us denote by $`A^U:T_pMT_pM`$ the symmetric operator $`g_M(A^U(X),Y)=g(_{}dF(X,Y),U)`$. From Lemma 2.1 and (2.3) we have
###### Proposition 2.1
If $`F`$ is an immersion without complex directions, then:
$`(i)`$ $`\mathrm{\Phi }`$ is parallel iff $`F^{}\omega `$ anti-commutes with $`A^U`$, $`UNM`$.
$`(ii)`$ If $`F`$ has equal Kähler angles, on $`\mathrm{\Omega }_{2n}^0`$, $`T^{}`$ is of type $`(1,1)`$ with respect to $`J_\omega `$.
$`(iii)`$ On $`\mathrm{\Omega }_{2n}^0`$, $`F`$ is pluriminimal iff $`T^{}`$ is of type $`(2,0)+(0,2)`$ with respect to $`J_\omega `$.
$`(iv)`$ If $`F`$ is broadly-pluriminimal, then, for $`p\mathrm{\Omega }_{2k}^0`$ with $`k1`$, $`T^{}`$ is of type $`(2,0)+(0,2)`$ with respect to any $`g_M`$-orthogonal complex structure $`\stackrel{~}{J}=J_\omega J^{}`$ on $`T_pM`$, where $`J^{}`$ is any $`g_M`$-orthogonal complex structure of $`𝒦_\omega `$.
Remark 1. If we call $`\omega _{NM}`$ the restriction of the Kähler form $`\omega `$ to the normal bundle $`NM`$, we see that, if $`\{X_\alpha ,Y_\alpha \}`$ is a diagonalizing $`g_M`$-orthonormal basis of $`F^{}\omega `$ at a point $`p`$, then $`\{U_\alpha =\mathrm{\Phi }(\frac{Y_\alpha }{\mathrm{sin}\theta _\alpha }),V_\alpha =\mathrm{\Phi }(\frac{X_\alpha }{\mathrm{sin}\theta _\alpha })\}`$ is a diagonalizing $`g`$-orthonormal basis of $`\omega _{NM}`$. Moreover, $`NM`$ has the same Kähler angles as $`F`$. Let $`J_{NM}`$ denote the complex structure on $`NM`$ defined by this basis, that is, the one that comes from the polar decomposition of $`\omega _{NM}`$. Then, $`\mathrm{\Phi }J_\omega =J_{NM}\mathrm{\Phi }`$.
We should also remark the following:
###### Proposition 2.2
If $`F`$ is an immersion with parallel 2-form $`F^{}\omega `$, then the Kähler angles are constant and, in particular, $`M=\mathrm{\Omega }_{2k}^0`$ for some $`k`$. In this case, considering $`TM`$ with the Levi-Civita connection $`_{}`$$`𝒦_\omega `$ and $`𝒦_\omega ^{}`$ are parallel sub-vector bundles of $`TM`$, and $`J_\omega C^{\mathrm{}}(𝒦_\omega ^{}𝒦_\omega ^{})`$, $`\stackrel{~}{g},\widehat{g}C^{\mathrm{}}(^2T^{}M)`$ are parallel sections. Furthermore, $`(X,Y,Z)g(_ZdF(X),JdF(Y))`$ is symmetric on $`TM`$, and, if $`F`$ has no complex directions, $`\widehat{}_{}=_{}`$. Moreover, if $`\mathrm{cos}\theta _{\alpha _1}>\mathrm{}>\mathrm{cos}\theta _{\alpha _r}`$ are the distinct eigenvalues of $`F^{}\omega `$, the corresponding eigenspaces $`E_{\alpha _t}`$ define a smooth integrable distribution of $`TM`$ whose integral submanifolds are parallel submanifolds of $`M`$. The integral submanifolds of $`E_{\alpha _r}`$ are isotropic in $`N`$ if $`\mathrm{cos}\theta _{\alpha _r}=0`$, and the ones of $`E_{\alpha _1}`$ are complex submanifolds of $`N`$ if $`\mathrm{cos}\theta _{\alpha _1}=1`$. The other ones are Kähler manifolds with respect to $`J_\omega `$, and $`F`$ restricted to each one of them is an immersion of constant equal Kähler angles $`\theta _{\alpha _t}`$ with respect to $`J`$.
Proof. If $`X,Y`$ are smooth vector fields on $`M`$ and $`ZT_pM`$, an elementary computation gives
$$_ZF^{}\omega (X,Y)=g(_ZdF(X),JdF(Y))+g(_ZdF(Y),JdF(X)),$$
(2.4)
which proves the symmetry of $`(X,Y,Z)g(_ZdF(X),JdF(Y))`$. From (2.2) we see that $`\widehat{g}`$ is parallel. Consequently, outside complex directions, $`_{}=\widehat{}_{}`$. If we parallel transport a diagonalizing orthonormal basis $`\{X_\alpha ,Y_\alpha \}`$ of $`F^{}\omega `$ at $`p_0`$ along geodesics, on a neighbourhood of $`p_0`$, since $`F^{}\omega `$ is parallel we get a diagonalizing orthonormal frame on a whole neighbourhood with the property $`_{}X_\alpha (p_0)=_{}Y_\alpha (p_0)=0`$. It also follows that $`\mathrm{cos}\theta _\alpha `$ remains constant along geodesics, so it is constant, and $`J_\omega (X_\alpha )=Y_\alpha `$ on a neighbourhood of $`p_0`$, with $`_{}J_\omega =0`$ at $`p_0`$, and so $`J_\omega `$ is parallel. Similarly we see that $`\stackrel{~}{g}`$ is parallel. If we extend $`F^{}\omega `$ to the complexified tangent space $`T_{p_0}^cM`$, then $`F^{}\omega (Z_\alpha )=i\mathrm{cos}\theta _\alpha Z_\alpha `$, and $`F^{}\omega (Z_{\overline{\alpha }})=i\mathrm{cos}\theta _\alpha Z_{\overline{\alpha }}`$. Obviouly, the corresponding eigenspaces of $`F^{}\omega `$, are parallel sub-vector bundles of $`T^cM`$. $`\mathbf{}`$
## 3 Immersions with equal Kähler angles
If $`F`$ has equal Kähler angles, then
$$F^{}\omega =\mathrm{cos}\theta J_\omega \text{and}\widehat{g}=\mathrm{sin}^2\theta g_M,$$
with $`\mathrm{cos}\theta `$ a locally Lipschitz map on $`M`$, smooth on the open set where it does not vanish, and $`\mathrm{\Omega }_{2k}^0=\mathrm{}`$ $`k0,n`$. Note that $`\mathrm{sin}^2\theta `$ and $`\mathrm{cos}^2\theta `$ are smooth on all $`M`$. The set $`=\mathrm{cos}\theta ^1(\{0\})`$ is the set of Lagrangian points, for, at these points, the tangent space of $`M`$ is a Lagrangian subspace of the tangent space of $`N`$. Its subset of interior points is $`\mathrm{\Omega }_0^0`$. Similarly, we say that $`𝒞=\mathrm{cos}\theta ^1(\{1\})`$ is the set of complex points. On the open set $`\mathrm{\Omega }_{2n}^0=\mathrm{cos}\theta ^1(IR\{0\})=M`$, $`J_\omega `$ defines a smooth almost complex structure $`g_M`$-orthogonal. On the open set $`\mathrm{cos}\theta ^1(IR\{1\})=M𝒞`$, $`\widehat{g}`$ is a smooth metric conformally equivalent to $`g_M`$. Thus, if $`n2`$, $`\widehat{}_{}=_{}`$ iff $`\theta `$ is constant. Since the Kähler angles are equal, any smooth local orthonormal frame of the type $`\{X_\alpha ,Y_\alpha =J_\omega X_\alpha \}`$ diagonalizes $`F^{}\omega `$ on the whole set where it is defined. From $`F^{}\omega =\mathrm{cos}\theta J_\omega `$, we get $`_XF^{}\omega =d\mathrm{cos}\theta (X)J_\omega +\mathrm{cos}\theta _XJ_\omega `$, with $`J_\omega `$ orthogonal to $`_XJ_\omega `$ with respect to the Hilbert-Schmidt inner product (because $`J_\omega ^2=2n`$ is constant). Hence, considering $`F^{}\omega `$ an operator on $`TM`$, on $`\mathrm{\Omega }_{2n}^0\mathrm{\Omega }_0^0`$
$$_{}F^{}\omega ^2=2n\mathrm{cos}\theta ^2+\mathrm{cos}^2\theta _{}J_\omega ^2.$$
(3.1)
We observe that $`M(\mathrm{\Omega }_{2n}^0\mathrm{\Omega }_0^0)`$ is a set of Lagrangian points with no interior. On $`\mathrm{\Omega }_{2n}^0`$, we have then, $`_{}F^{}\omega =0`$ iff $`_{}J_\omega =0`$ and $`\theta `$ is constant. Note that $`_{}F^{}\omega ^2`$, considering $`F^{}\omega `$ an operator on $`TM`$, is twice the square norm when considering $`F^{}\omega `$ a 2-form. From (2.3) we get, on $`M𝒞`$,
$$\mathrm{\Phi }(T^{}(X,Y))=2\mathrm{cos}\theta (_{}dF)^{(1,1)}(J_\omega X,Y).$$
(3.2)
The right-hand side of (3.2) is defined to be zero at a Lagrangian point. Consequentely
###### Proposition 3.1
If $`F`$ is an immersion with equal Kähler angles and without complex points, then $`T^{}=0`$, that is, $`_{}^{}=\widehat{}_{}`$ iff $`\mathrm{\Phi }`$ is parallel iff $`F`$ is Lagrangian or pluriminimal. In particular, if $`F`$ is minimal, $`\mathrm{\Phi }`$ is parallel iff $`F`$ is broadly-pluriminimal.
Let $`Re(u+iv)=u`$, for $`u,vNM`$.
###### Proposition 3.2
If $`F`$ is any immersion with equal Kähler angles, then, outside complex and Lagrangian points,
$$\mathrm{\Phi }\left(\frac{1n}{4}\mathrm{log}\mathrm{sin}^2\theta \right)=\frac{4\mathrm{cos}\theta }{\mathrm{sin}^2\theta }Re\left(i\underset{\beta ,\mu }{}\left(g(_{\overline{\mu }}dF(\mu ),JdF(\beta ))g(_{\overline{\mu }}dF(\beta ),JdF(\mu ))\right)\mathrm{\Phi }(\overline{\beta })\right),$$
where $`\mathrm{log}\mathrm{sin}^2\theta `$ is the gradient with respect to $`g_M`$.
If $`F`$ is a complex submanifold on a open set, then $`J_\omega `$ is the induced complex structure on $`M`$ and $`_{}dF`$ is of type $`(2,0)`$. Applying Proposition 2.2 on $`\mathrm{\Omega }_0^0`$, and Proposition 3.1 on open sets without complex and Lagrangian points, and noting that $`\{\mathrm{\Phi }(\beta ),\mathrm{\Phi }(\overline{\beta })=\overline{\mathrm{\Phi }(\beta )}\}_{1\beta n}`$ multiplied by $`\frac{\sqrt{2}}{\mathrm{sin}\theta }`$ constitutes an unitary basis of $`NM^c`$, we immediately conclude
###### Corollary 3.1
If $`F`$ is an immersion with equal Kähler angles, and $`n2`$, then $`\theta `$ is constant iff
$$\underset{\mu }{}g(_{\overline{\mu }}dF(\mu ),JdF(\beta ))=\underset{\mu }{}g(_{\overline{\mu }}dF(\beta ),JdF(\mu ))\beta .$$
(3.3)
Note that (3.3) is a sort of symmetry property, and the first term is just $`\frac{n}{2}g(H,JdF(\beta ))`$, where $`H=\frac{1}{2n}trace_{g_M}_{}dF=\frac{2}{n}_\mu _{}dF(\overline{\mu },\mu )`$ is the mean curvature of $`F`$.
###### Theorem 3.1
If $`n2`$ and $`F`$ is a pluriminimal immersion with equal Kähler angles then $`\mathrm{cos}\theta =`$ constant. Moreover, if it is not a complex submanifold, then $`_{}=\widehat{}_{}=_{}^{}`$, and $`N`$ must be Ricci-flat. In particular, $`\mathrm{\Phi }`$ defines a parallel homothetic isomorphism from $`(TM,g_M)`$ onto $`(NM,g)`$.
Proof. On a neighbourhood of a non-complex point, from Proposition 3.1, $`\widehat{}_{}=_{}^{}`$, and from Corollary 3.1, $`\mathrm{cos}\theta `$ is constant. Then $`\widehat{}_{}=_{}`$, as well. So if $`F`$ is not a complex submanifold, it has no complex points anywhere. Finally, (1.7) for pluriminimal immersions with $`\kappa =constant`$ gives $`R=0`$. $`\mathbf{}`$
The above theorem and Proposition 3.1 lead to:
###### Corollary 3.2
If $`F`$ is a minimal immersion with equal Kähler angles, without complex points, $`n2`$, and $`R0`$, then $`F`$ is Lagrangian iff $`\mathrm{\Phi }`$ is parallel.
To prove Proposition 3.2 we will need to relate the three connections of $`M`$, $`_{}`$, $`\widehat{}_{}`$ and $`_{}^{}`$. Let $`\{e_1,\mathrm{},e_{2n}\}=\{X_\mu ,Y_\mu =J_\omega X_\mu \}_{1\mu n}`$ be a local $`g_M`$-orthonormal frame outside the Lagrangian and complex set, and $`_1,\mathrm{},_{2n}`$ a local frame of $`M`$ defined by a coordinate chart. Set $`g_{ij}=g_M(_i,_j)`$$`\widehat{g}_{ij}=\widehat{g}(_i,_j)=\mathrm{sin}^2\theta g_{ij}`$, and $`e_s=_i\lambda _{si}_i`$. The Christofel symbols are given by $`2\widehat{\mathrm{\Gamma }}_{ij}^k=_s\widehat{g}^{ks}(_i\widehat{g}_{sj}+_j\widehat{g}_{is}_s\widehat{g}_{ij})=\delta _{kj}_i\mathrm{log}\mathrm{sin}^2\theta +\delta _{ki}_j\mathrm{log}\mathrm{sin}^2\theta _sg^{ks}g_{ij}_s\mathrm{log}\mathrm{sin}^2\theta +2\mathrm{\Gamma }_{ij}^k.`$ Hence
$$\widehat{}__i_j__i_j=\underset{k}{}(\widehat{\mathrm{\Gamma }}_{ij}^k\mathrm{\Gamma }_{ij}^k)_k=\frac{1}{2}(_i(\mathrm{log}\mathrm{sin}^2\theta )_j+_j(\mathrm{log}\mathrm{sin}^2\theta )_ig_{ij}(\mathrm{log}\mathrm{sin}^2\theta ))\text{}$$
Since $`_{ij}g_{ij}\lambda _{si}\lambda _{sj}=1`$, $`_s\widehat{}_{e_s}e_s_{e_s}e_s=_{sij}\lambda _{si}\lambda _{sj}(\widehat{}__i_j__i_j)=(1n)\mathrm{log}\mathrm{sin}^2\theta `$. Therefore,
$`{\displaystyle \underset{\mu }{}}\widehat{}_{\overline{\mu }}\mu _{\overline{\mu }}\mu ={\displaystyle \frac{1}{4}}{\displaystyle \underset{\mu }{}}\left(\widehat{}_{X_\mu }X_\mu +\widehat{}_{Y_\mu }Y_\mu _{X_\mu }X_\mu _{Y_\mu }Y_\mu \right)i\left(\widehat{}_{X_\mu }Y_\mu \widehat{}_{Y_\mu }X_\mu _{X_\mu }Y_\mu +_{Y_\mu }X_\mu \right)`$ (3.4)
$`={\displaystyle \frac{1}{4}}{\displaystyle \underset{s}{}}(\widehat{}_{e_s}e_s_{e_s}e_s)+{\displaystyle \frac{i}{4}}{\displaystyle \underset{\mu }{}}\left([Y_\mu ,X_\mu ][Y_\mu ,X_\mu ]\right)={\displaystyle \frac{(1n)}{4}}\mathrm{log}\mathrm{sin}^2\theta .`$
Set $`S^{}(X,Y)=_X^{}Y\widehat{}_XY`$. Then $`S^{}(X,Y)S^{}(Y,X)=T^{}(X,Y)`$. Similarly we get
$$\underset{\mu }{}_{\overline{\mu }}^{}\mu \widehat{}_{\overline{\mu }}\mu =\frac{1}{4}trace_{g_M}S^{}\frac{i}{4}\underset{\mu }{}T^{}(X_\mu ,Y_\mu ).$$
(3.5)
###### Lemma 3.1
$`XT_pM`$, $`_i\widehat{g}(S^{}(e_i,e_i),X)=_i\widehat{g}(T^{}(e_i,X),e_i)`$.
Proof. We may assume that the local referencial $`_i`$ is $`\widehat{g}`$-orthonormal at a fixed poit $`p_0`$. On a neighbourhood of $`p_0`$, we define $`\mathrm{\Gamma }_{}^{}{}_{ij}{}^{k}`$ and $`S_{}^{}{}_{ij}{}^{k}`$ as
$$__i^{}_j=\underset{k}{}\mathrm{\Gamma }_{}^{}{}_{ij}{}^{k}_kS^{}(_i,_j)=\underset{k}{}S_{}^{}{}_{ij}{}^{k}_k=\underset{k}{}(\mathrm{\Gamma }_{}^{}{}_{ij}{}^{k}\widehat{\mathrm{\Gamma }}_{ij}^k)_k.$$
Then $`T_{ij}^{}{}_{}{}^{k}=\mathrm{\Gamma }_{}^{}{}_{ij}{}^{k}\mathrm{\Gamma }_{}^{}{}_{ji}{}^{k}`$, and, at $`p_0`$$`\mathrm{\Gamma }_{}^{}{}_{ij}{}^{k}=\widehat{g}(__i^{}_j,_k),S_{}^{}{}_{ij}{}^{k}=\widehat{g}(S^{}(_i,_j),_k)=\mathrm{\Gamma }_{}^{}{}_{ij}{}^{k}\widehat{\mathrm{\Gamma }}_{ij}^k`$. $`_{}^{}`$ is a Riemannian connection with respect to $`\widehat{g}`$. Then
$$_i\widehat{g}_{jk}(p_0)=\widehat{g}(__i^{}_j,_k)+\widehat{g}(_j,__i^{}_k)=\mathrm{\Gamma }_{}^{}{}_{ij}{}^{k}+\mathrm{\Gamma }_{}^{}{}_{ik}{}^{j}$$
Hence, at $`p_0`$
$`2\widehat{\mathrm{\Gamma }}_{ij}^k`$ $`=`$ $`{\displaystyle \underset{s}{}}\widehat{g}^{ks}(_i\widehat{g}_{sj}+_j\widehat{g}_{is}_s\widehat{g}_{ij})=\mathrm{\Gamma }_{}^{}{}_{ik}{}^{j}+\mathrm{\Gamma }_{}^{}{}_{ij}{}^{k}+\mathrm{\Gamma }_{}^{}{}_{ji}{}^{k}+\mathrm{\Gamma }_{}^{}{}_{jk}{}^{i}\mathrm{\Gamma }_{}^{}{}_{ki}{}^{j}\mathrm{\Gamma }_{}^{}{}_{kj}{}^{i}`$
$`=`$ $`(\mathrm{\Gamma }_{}^{}{}_{ij}{}^{k}+\mathrm{\Gamma }_{}^{}{}_{ji}{}^{k})+(\mathrm{\Gamma }_{}^{}{}_{ik}{}^{j}\mathrm{\Gamma }_{}^{}{}_{ki}{}^{j})+(\mathrm{\Gamma }_{}^{}{}_{jk}{}^{i}\mathrm{\Gamma }_{}^{}{}_{kj}{}^{i})=(\mathrm{\Gamma }_{}^{}{}_{ij}{}^{k}+\mathrm{\Gamma }_{}^{}{}_{ji}{}^{k})+T_{}^{}{}_{ik}{}^{j}+T_{}^{}{}_{jk}{}^{i}`$
But $`\mathrm{\Gamma }_{}^{}{}_{ij}{}^{k}+\mathrm{\Gamma }_{}^{}{}_{ji}{}^{k}=2\mathrm{\Gamma }_{}^{}{}_{ij}{}^{k}+(\mathrm{\Gamma }_{}^{}{}_{ji}{}^{k}\mathrm{\Gamma }_{}^{}{}_{ij}{}^{k})=2\mathrm{\Gamma }_{}^{}{}_{ij}{}^{k}+T_{}^{}{}_{ji}{}^{k}`$. Thus
$$S_{}^{}{}_{ij}{}^{k}=\mathrm{\Gamma }_{}^{}{}_{ij}{}^{k}\widehat{\mathrm{\Gamma }}_{ij}^k=\frac{1}{2}(T_{}^{}{}_{ij}{}^{k}T_{}^{}{}_{ik}{}^{j}+T_{}^{}{}_{kj}{}^{i}).$$
That is, at $`p_0`$$`\widehat{g}(S^{}(_i,_j),_k)=\frac{1}{2}\left(\widehat{g}(T^{}(_i,_j),_k)\widehat{g}(T^{}(_i,_k),_j)+\widehat{g}(T^{}(_k,_j),_i)\right).`$ We may assume that, at $`p_0`$, $`_i(p_0)=\frac{e_i}{\mathrm{sin}\theta }`$, leading to the Lemma. $`\mathbf{}`$
Proof of Proposition 3.2. Following the proof of Lemma 2.1, $`\mathrm{\Phi }(_X^{}\mu _X\mu )=`$
$`=((Ji\mathrm{cos}\theta )_XdF(\mu ))^{}`$. Hence, from (3.4),
$$\mathrm{\Phi }(\frac{(1n)}{4}\mathrm{log}\mathrm{sin}^2\theta )=\mathrm{\Phi }(\underset{\mu }{}\widehat{}_{\overline{\mu }}\mu _{\overline{\mu }}\mu )=((Ji\mathrm{cos}\theta )\frac{nH}{2})^{}\underset{\mu }{}\mathrm{\Phi }(_{\overline{\mu }}^{}\mu \widehat{}_{\overline{\mu }}\mu ).$$
But, from (3.5), $`_\mu \mathrm{\Phi }(_{\overline{\mu }}^{}\mu \widehat{}_{\overline{\mu }}\mu )=\frac{1}{4}\mathrm{\Phi }(trace_{g_M}S^{})\frac{i}{4}\mathrm{\Phi }(_\mu T^{}(X_\mu ,Y_\mu ))`$. The skew-symmetry of $`T^{}`$ and (3.2) implies that
$$\mathrm{\Phi }(\underset{\mu }{}T^{}(X_\mu ,Y_\mu ))=2i\underset{\mu }{}\mathrm{\Phi }(T^{}(\mu ,\overline{\mu }))=4\mathrm{cos}\theta _\mu dF(\overline{\mu })=2n\mathrm{cos}\theta H.$$
Thus, $`_\mu \mathrm{\Phi }(_{\overline{\mu }}^{}\mu \widehat{}_{\overline{\mu }}\mu )=\frac{1}{4}\mathrm{\Phi }(trace_{g_M}S^{})\frac{ni}{2}\mathrm{cos}\theta H.`$ Therefore,
$$\mathrm{\Phi }(\frac{(1n)}{4}\mathrm{log}\mathrm{sin}^2\theta )=\frac{1}{4}\left(2n(JH)^{}\mathrm{\Phi }(Trace_{g_M}S^{})\right).$$
(3.6)
Using Lemma 3.1, (3.2), and $`\mathrm{\Phi }(\mu )=JdF(\mu )i\mathrm{cos}\theta dF(\mu )`$, we have
$`\mathrm{\Phi }(Trace_{g_M}S^{})={\displaystyle \underset{j,k}{}}\widehat{g}(S^{}(e_j,e_j),{\displaystyle \frac{e_k}{\mathrm{sin}\theta }})\mathrm{\Phi }({\displaystyle \frac{e_k}{\mathrm{sin}\theta }})={\displaystyle \underset{j,k}{}}\widehat{g}(T^{}(e_j,{\displaystyle \frac{e_k}{\mathrm{sin}\theta }}),e_j)\mathrm{\Phi }({\displaystyle \frac{e_k}{\mathrm{sin}\theta }})`$
$`={\displaystyle \frac{4}{\mathrm{sin}^2\theta }}{\displaystyle \underset{\mu ,\beta }{}}(\left(\widehat{g}(T^{}(\mu ,\beta ),\overline{\mu })+\widehat{g}(T^{}(\overline{\mu },\beta ),\mu )\right)\mathrm{\Phi }(\overline{\beta })+\left(\widehat{g}(T^{}(\mu ,\overline{\beta }),\overline{\mu })+\widehat{g}(T^{}(\overline{\mu },\overline{\beta }),\mu )\right)\mathrm{\Phi }(\beta ))`$
$`={\displaystyle \frac{4}{\mathrm{sin}^2\theta }}{\displaystyle \underset{\mu ,\beta }{}}\left(g(\mathrm{\Phi }\left(T^{}(\overline{\mu },\beta )\right),\mathrm{\Phi }(\mu ))\mathrm{\Phi }(\overline{\beta })+g(\mathrm{\Phi }\left(T^{}(\mu ,\overline{\beta })\right),\mathrm{\Phi }(\overline{\mu }))\mathrm{\Phi }(\beta )\right)`$
$`={\displaystyle \frac{8i\mathrm{cos}\theta }{\mathrm{sin}^2\theta }}{\displaystyle \underset{\mu ,\beta }{}}\left(g(_{\overline{\mu }}dF(\beta ),JdF(\mu ))\mathrm{\Phi }(\overline{\beta })g(_\mu dF(\overline{\beta }),JdF(\overline{\mu }))\mathrm{\Phi }(\beta )\right).`$
Writing $`(JH)^{}`$ in terms of $`\mathrm{\Phi }(\beta )`$ and $`\mathrm{\Phi }(\overline{\beta })`$,
$`2n(JH)^{}`$ $`=`$ $`{\displaystyle \underset{\beta }{}}{\displaystyle \frac{4n}{\mathrm{sin}^2\theta }}\left(g(JH,\mathrm{\Phi }(\beta ))\mathrm{\Phi }(\overline{\beta })+g(JH,\mathrm{\Phi }(\overline{\beta }))\mathrm{\Phi }(\beta )\right)`$
$`=`$ $`{\displaystyle \underset{\beta ,\mu }{}}{\displaystyle \frac{8i\mathrm{cos}\theta }{\mathrm{sin}^2\theta }}\left(g(_{\overline{\mu }}dF(\mu ),JdF(\beta ))\mathrm{\Phi }(\overline{\beta })g(_{\overline{\mu }}dF(\mu ),JdF(\overline{\beta }))\mathrm{\Phi }(\beta )\right),`$
and substituing these equations into (3.6), we prove Proposition 3.2. $`\mathbf{}`$
### 3.1 The Weitzenböck formula for $`F^{}\omega `$
For simplicity let us use the notation
$$g_XYZ=g(_XdF(Y),JdF(Z)).$$
We also observe that, from
$$\mu \frac{i}{2}\mathrm{cos}\theta =F^{}\omega (\mu ,\overline{\mu }),$$
(3.7)
valid on an open set, and from (2.4), we obtain $`\mu `$
$`{\displaystyle \frac{i}{2}}d\mathrm{cos}\theta (X)`$ $`=`$ $`d(F^{}\omega (\mu ,\overline{\mu }))(X)=_XF^{}\omega (\mu ,\overline{\mu })+F^{}\omega (_X\mu ,\overline{\mu })+F^{}\omega (\mu ,_X\overline{\mu })`$ (3.8)
$`=`$ $`g_X\mu \overline{\mu }+g_X\overline{\mu }\mu +2(_X\mu ,\overline{\mu }+_X\overline{\mu },\mu )F^{}\omega (\mu ,\overline{\mu })`$
$`=`$ $`g_X\mu \overline{\mu }+g_X\overline{\mu }\mu \text{ (no sumation over }\mu ).`$
Then (3.3) is equivalent to $`g(_XdF(\mu ),JdF(\overline{\mu }))=g(_XdF(\overline{\mu }),JdF(\mu )),\mu `$ (or some $`\mu `$). From $`J_\omega Z_\alpha =iZ_\alpha `$, $`J_\omega Z_{\overline{\alpha }}=iZ_{\overline{\alpha }}`$ and the fact that $`J_\omega `$ is $`g_M`$-orthogonal, we get, on $`\mathrm{\Omega }_{2n}^0`$, $`\alpha ,\beta `$, and $`vTM`$
$$_vJ_\omega (\alpha ),\beta =2i_v\alpha ,\beta ,_vJ_\omega (\alpha ),\overline{\beta }=0.$$
(3.9)
Recall that, if $`\xi `$ is a $`r+1`$-form on $`M`$, $`r0`$, with values on a vector bundle $`E`$ over $`M`$ with a connection $`_{}^E`$, then $`\delta \xi `$, the divergence of $`\xi `$, is the $`r`$-form on $`M`$ with values on $`E`$ given by
$$\delta \xi (u_1,\mathrm{},u_r)=\underset{s}{}_{e_s}^E\xi (e_s,u_1,\mathrm{},u_r),$$
where $`e_1,\mathrm{},e_m`$ is an orthonormal basis of $`T_pM`$, $`u_iT_pM`$, and $`_{}^E\xi `$ is the covariant derivative of $`\xi `$ on $`^{r+1}T^{}ME`$. Thus, $`\delta `$ is the formal adjoint of $`d`$ on forms (cf. \[E-L\]). Note that $`\delta F^{}\omega (X)=\delta F^{}\omega ,X,XT_pM`$, considering on the left-hand side $`F^{}\omega `$ a (closed) 2-form on $`M`$ and on the right-hand side an endomorphism of $`TM`$.
###### Proposition 3.3
Let $`F`$ be an immersion with equal Kähler angles and $`\mathrm{cos}\theta `$ denote the gradient with respect to $`g_M`$. On $`\mathrm{\Omega }_{2n}^0`$, and considering $`F^{}\omega `$ an endomorphism of $`TM`$.
$$\delta F^{}\omega =(n2)J_\omega (\mathrm{cos}\theta ),\mathrm{cos}\theta (\delta J_\omega )=(n1)J_\omega (\mathrm{cos}\theta ).$$
Thus,
$`(i)`$ For $`n=1`$, $`\delta J_\omega =0`$ (obviously!), and $`\delta F^{}\omega =0`$ iff $`\theta `$ is constant.
$`(ii)`$ For $`n=2`$, $`\delta F^{}\omega =0`$. Moreover, $`\delta J_\omega =0`$ iff $`\theta `$ is constant.
$`(iii)`$ For $`n1,2`$$`\delta F^{}\omega =0`$ iff $`\delta J_\omega =0`$ iff $`\theta `$ is constant.
Proof. Considering $`F^{}\omega `$ a 2-form on $`M`$, using the symmetry of $`_{}dF`$ and (2.4), if $`XT_pM`$,
$`\delta (F^{}\omega )(X)=`$ $`{\displaystyle \underset{\mu }{}}2_\mu F^{}\omega (\overline{\mu },X)2_{\overline{\mu }}F^{}\omega (\mu ,X)={\displaystyle \underset{\mu }{}}2g_\mu \overline{\mu }X2g_\mu X\overline{\mu }+2g_{\overline{\mu }}\mu X2g_{\overline{\mu }}X\mu `$
$`=`$ $`2{\displaystyle \underset{\mu }{}}(g_X\mu \overline{\mu }+g_X\overline{\mu }\mu )4{\displaystyle \underset{\mu }{}}(g_{\overline{\mu }}X\mu g_{\overline{\mu }}\mu X).`$
From (3.8), $`\frac{ni}{2}d\mathrm{cos}\theta (X)=_\mu g_X\mu \overline{\mu }+g_X\overline{\mu }\mu `$. Therefore,
$$\delta (F^{}\omega )(X)=nid\mathrm{cos}\theta (X)4\underset{\mu }{}_{\overline{\mu }}F^{}\omega (\mu ,X).$$
(3.10)
Since $`F^{}\omega `$ is of type $`(1,1)`$ with respect to $`J_\omega `$, and $`ZT_p^cM`$, $`\mu ,\beta `$, $`_Z\beta ,\mu =\beta ,_Z\mu `$, we get using (3.9)
$`_ZF^{}\omega (\mu ,\beta )`$ $`=`$ $`d(F^{}\omega (\mu ,\beta ))(Z)F^{}\omega (_Z\mu ,\beta )F^{}\omega (\mu ,_Z\beta )`$ (3.11)
$`=`$ $`2i\mathrm{cos}\theta _Z\mu ,\beta =\mathrm{cos}\theta _ZJ_\omega (\mu ),\beta .`$
Note that, since $`J_\omega ^2=Id`$, $`_XJ_\omega (J_\omega Y)=J_\omega (_XJ_\omega (Y))`$, $`X,YT_pM`$. So
$`4{\displaystyle \underset{\mu }{}}_{\overline{\mu }}J_\omega (\mu )`$ $`=`$ $`{\displaystyle \underset{\mu }{}}_{X_\mu }J_\omega (X_\mu )+_{Y_\mu }J_\omega (Y_\mu )+i_{Y_\mu }J_\omega (X_\mu )i_{X_\mu }J_\omega (Y_\mu )`$
$`=`$ $`\delta J_\omega +i{\displaystyle \underset{\mu }{}}(_{X_\mu }J_\omega (J_\omega X_\mu )_{Y_\mu }J_\omega (J_\omega Y_\mu ))=\left(\delta J_\omega +iJ_\omega (\delta J_\omega )\right).`$
Hence, from (3.11), and since $`J_\omega `$ is $`g_M`$-orthogonal, $`\beta `$
$$\underset{\mu }{}_{\overline{\mu }}F^{}\omega (\mu ,\beta )=\frac{\mathrm{cos}\theta }{4}\delta J_\omega +iJ_\omega (\delta J_\omega ),\beta =\frac{\mathrm{cos}\theta }{2}\delta J_\omega ,\beta .$$
Moreover, $`id\mathrm{cos}\theta (\beta )=d\mathrm{cos}\theta (J_\omega \beta )=\mathrm{cos}\theta ,J_\omega \beta =J_\omega (\mathrm{cos}\theta ),\beta .`$ From (3.10), $`\delta F^{}\omega (\beta )=nJ_\omega (\mathrm{cos}\theta )+2\mathrm{cos}\theta \delta J_\omega ,\beta `$. Thus, if we consider $`F^{}\omega `$ an endomorphism of $`TM`$, and since $`,`$, $`J_\omega `$, and $`F^{}\omega `$ are real operators,
$$\delta F^{}\omega =nJ_\omega (\mathrm{cos}\theta )+2\mathrm{cos}\theta \delta J_\omega .$$
(3.12)
On the other hand, $`F^{}\omega =\mathrm{cos}\theta J_\omega `$. Then, an elementary computation gives
$$\delta F^{}\omega =J_\omega (\mathrm{cos}\theta )+\mathrm{cos}\theta \delta J_\omega .$$
(3.13)
Comparing (3.12) with (3.13) we get the Proposition. $`\mathbf{}`$
Remark 2. One may check the equation in Proposition 3.2 by using the equalities given in the above Proposition and its proof.
If we apply the Weitzenböck formula to the 2-form $`F^{}\omega `$, for an immersion $`F:MN`$ we get (see e.g \[E-L\] (1.32))
$$\frac{1}{2}\mathrm{}F^{}\omega ^2=\mathrm{}F^{}\omega ,F^{}\omega +_{}F^{}\omega ^2+SF^{}\omega ,F^{}\omega ,$$
(3.14)
where $`,`$ denotes the Hilbert-Schmidt inner product for 2-forms, and $`S`$ is the Ricci operator of $`^2T^{}M`$. We note that we use the the sign convention $`\mathrm{}\varphi =+Trace_{g_M}Hess\varphi `$, for $`\varphi `$ a smooth real map on $`M`$. This sign is opposite to the one of \[E-L\], but here we use the same sign as in \[E-L\] for the Laplacian of forms $`\mathrm{}=d\delta +\delta d`$. If $`\overline{R}`$ denotes the curvature tensor of $`^2T^{}M`$, and $`X,Y,u,vT_pM`$, $`\xi ^2T_p^{}M`$, then
$`\overline{R}(X,Y)\xi (u,v)=\xi (R^M(X,Y)u,v)\xi (u,R^M(X,Y)v),`$
$`SF^{}\omega (X,Y)={\displaystyle \underset{1i2n}{}}\overline{R}(e_i,X)F^{}\omega (e_i,Y)+\overline{R}(e_i,Y)F^{}\omega (e_i,X),`$
Where $`R^M`$ denotes the curvature tensor of $`M`$. In general, we use the following sign convention for curvature tensors: $`R^M(X,Y)Z=_X_YZ+_Y_XZ+_{[X,Y]}Z`$. Then, $`R^M(X,Y,Z,W)=g_M(R^M(X,Y)Z,W)`$. It is straightforward to prove
###### Lemma 3.2
If $`\{X_\alpha ,Y_\alpha \}`$ is a diagonalizing orthonormal basis of $`F^{}\omega `$ at $`p`$,
$`SF^{}\omega ,F^{}\omega `$ $`=`$ $`{\displaystyle \underset{\mu }{}}4\mathrm{cos}^2\theta _\mu Ricci^M(\mu ,\overline{\mu })+{\displaystyle \underset{\mu ,\rho }{}}8\mathrm{cos}\theta _\mu \mathrm{cos}\theta _\rho R^M(\rho ,\overline{\rho },\mu ,\overline{\mu })`$
$`=`$ $`{\displaystyle \underset{\mu ,\rho }{}}4(\mathrm{cos}\theta _\mu +\mathrm{cos}\theta _\rho )^2R^M(\rho ,\mu ,\overline{\rho },\overline{\mu })+4(\mathrm{cos}\theta _\mu \mathrm{cos}\theta _\rho )^2R^M(\overline{\rho },\mu ,\rho ,\overline{\mu }).`$
In particular, if $`F`$ has equal Kähler angles at $`p`$, then, at $`p`$,
$$SF^{}\omega ,F^{}\omega =16\mathrm{cos}^2\theta \underset{\rho ,\mu }{}R^M(\rho ,\mu ,\overline{\rho },\overline{\mu }).$$
Moreover, if $`(M,J_\omega ,g_M)`$ is Kähler in a neighbourhood of $`p`$, then $`SF^{}\omega ,F^{}\omega =0`$.
For example, if $`M`$ has constant sectional curvature $`K`$, $`SF^{}\omega ,F^{}\omega =4(n1)KF^{}\omega ^2`$. If $`(M,J_\omega ,g_M)`$ is a Kähler manifold of constant holomorphic sectional curvature $`K`$ then $`SF^{}\omega ,F^{}\omega `$ $`=4K\left(n_\mu \mathrm{cos}^2\theta _\mu \left(_\mu \mathrm{cos}\theta _\mu \right)^2\right)`$ has constant sign, with equality to zero iff $`K=0`$ or $`F`$ has equal Kähler angles. If $`F^{}\omega `$ is parallel, from (3.14), we obtain that $`SF^{}\omega ,F^{}\omega =0`$. In the latter case, if $`n2`$ and $`M`$ has constant sectional curvature, then, either $`F`$ is Lagrangian, or $`K=0`$.
We recall the concept of non-negative isotropic sectional curvature, for $`M`$ with dimension $`4`$, defined by Micallef and Moore in \[Mi-Mo\]. Let
$$K_{isot}(\sigma )=\frac{R^M(z,w,\overline{z},\overline{w})}{zw^2},$$
where $`\sigma =span_{lC}\{z,w\}`$ is a totally isotropic complex two-plane in $`T^cM`$, that is, $`u\sigma g_M(u,u)=0`$, and where $`R^M(x,y,u,v)`$ is extend to the complexified tangent space by $`lC`$-multilinearity. The curvature of $`M`$ is said to be non-negative (resp. positive) on totally isotropic two-planes at $`p`$, if $`K(\sigma )0`$ (resp. $`>0`$) whenever $`\sigma T_p^cM`$ is a totally isotropic two-plane over $`p`$. If $`M`$ is compact, simply connected and has positive isotropic sectional curvature everywhere, then $`M`$ is homeomorphic to a sphere (\[Mi-Mo\]). If $`n1`$, $`T^{2n}`$ is the flat torus, and $`S^2`$ is the euclidean sphere of $`IR^3`$, then $`S^2\times T^{2n}`$, $`S^2\times S^2`$, $`S^2\times S^2\times T^{2n}`$ have isotropic sectional curvature $`0`$ but not $`>0`$. If $`\{X_\alpha ,Y_\alpha \}`$ is any orthonormal basis of $`T_pM`$, and $`\mathrm{`}\mathrm{`}\mu \mathrm{"}`$ denotes $`Z_\mu `$ as in (1.3), the expression
$$S_{isot}(\{Z_\alpha \}_{1\alpha n})=\underset{\rho \mu }{}K_{isot}(span_{lC}\{\rho ,\mu \})=4\underset{\rho ,\mu }{}R^M(\rho ,\mu ,\overline{\rho },\overline{\mu })$$
(3.15)
is a hermitian trace of the curvature of $`M`$ restricted to the maximal totally isotropic subspace $`span_{lC}\{Z_1,\mathrm{},Z_n\}`$ of $`T^cM`$. To require it to be $`0`$, for all maximal totally isotropic subspaces - and we will say that $`M`$ has non-negative isotropic scalar curvature \- seems to be strictly weaker than to have non-negative isotropic sectional curvature. We also note that, any other metric conformaly equivalent to the flat metric $`g_0`$ on the 2n-torus with non-negative isotropic scalar curvature is homothetically equivalent to $`g_0`$, hence flat. In fact, in general, if $`\widehat{g}=e^\varphi g_M`$ is a conformaly equivalent metric on $`M`$, then, for each $`g_M`$-orthonormal basis $`\{X_\alpha ,Y_\alpha \}`$, $`\widehat{S}_{isot}(\{\widehat{Z}_\alpha \})=e^\varphi S_{isot}(\{Z_\alpha \})(n1)e^{2\varphi }(2\mathrm{}\varphi +(n1)\varphi ^2)`$, where $`\widehat{Z_\alpha }`$ are defined by the $`\widehat{g}`$-orthonormal basis $`\{e^{\frac{\varphi }{2}}X_\alpha ,e^{\frac{\varphi }{2}}Y_\alpha \}`$. To require $`2\mathrm{}\varphi +(n1)\varphi ^20`$, implies, in case of $`M`$ compact, $`\varphi `$ constant. We observe that, if $`dim_{IR}M6`$, then $`S_{isot}0`$ does not imply $`M`$ to be flat, but $`K_{isot}0`$ implies so. We also note that, if $`dim_{IR}(T_pM)=4`$, the set of curvature operators at $`p`$ with zero isotropic sectional curvature, is a vector space of dimension 9.
Recall that, for an immersion with equal Kähler angles, $`F^{}\omega `$ is parallel iff $`\theta `$ is constant and if $`\mathrm{cos}\theta 0`$, $`(M,J_\omega ,g_M)`$ is a Kähler manifold. We are going to see that an extra condition on the scalar isotropic curvature of $`M`$ may imply itself that the Kähler angle is constant and/or $`_{}J_\omega =0`$. From Proposition 3.3, for any $`n1`$, on $`\mathrm{\Omega }_{2n}^0\mathrm{\Omega }_0^0`$
$$\delta F^{}\omega ^2=(n2)^2\mathrm{cos}\theta ^2.$$
(3.16)
In particular, if $`n2`$, $`\mathrm{cos}\theta ^2`$ can be extended as a smooth map on all $`M`$ (recall that $`\mathrm{\Omega }_{2n}^0\mathrm{\Omega }_0^0`$ is dense on $`M`$), and from (3.1) we get that $`\mathrm{cos}^2\theta _{}J_\omega ^2`$ is also smooth. Observe that $`\delta F^{}\omega ^2`$ has the same value considering $`\delta F^{}\omega `$ a vector or a 1-form, but considering $`F^{}\omega `$ a 2-form (as in (3.14)) $`_{}F^{}\omega ^2`$ is half of the square norm when considering $`F^{}\omega `$ an operator of $`TM`$ (as in (3.1)). For $`n=2`$, $`F^{}\omega `$ is co-closed, and so it is a harmonic 2-form. In fact, since $`F`$ has equal Kähler angles, $`F^{}\omega =\mathrm{cos}\theta (X_{}^1Y_{}^1+X_{}^2Y_{}^2)`$, and so $`F^{}\omega =\pm F^{}\omega `$, where $``$ is the Hodge star-operator of $`(M,g)`$, and the $`\pm `$ sign depends on the orientation of the diagonalizing basis. In particular, $`F^{}\omega `$ is co-closed. For $`n3`$, $`F^{}\omega `$ is harmonic iff $`\theta `$ is constant.
Integrating (3.14) on $`M`$, using (3.16) and (3.1), and the fact that $`_M\mathrm{}F^{}\omega ,F^{}\omega Vol_M=_M\delta F^{}\omega ^2Vol_M`$, we have
$$0=_M\left((n(n2)^2)\mathrm{cos}\theta ^2+\frac{1}{2}\mathrm{cos}^2\theta _{}J_\omega ^2\right)Vol_M+_MSF^{}\omega ,F^{}\omega Vol_M.$$
(3.17)
The first integrand is smooth on $`M`$, for all $`n`$ ( for $`n=2`$ it gives half of (3.1)). The factor $`n(n2)^2`$ is respectively, $`>0`$, $`=0`$, $`<0`$, according $`n=2`$ or $`3`$, $`n=4`$, and $`n5`$. If $`M`$ has non-negative isotropic scalar curvature, $`SF^{}\omega ,F^{}\omega 0`$, by Lemma 3.2. We conclude:
###### Proposition 3.4
Let $`F`$ be a non-Lagrangian immersion with equal Kähler angles of a compact orientable $`M`$ with non-negative isotropic scalar curvature into a Kähler manifold $`N`$. If $`n=2`$ or $`3`$, then $`\theta `$ is constant and $`(M,J_\omega ,g_M)`$ is a Kähler manifold. If $`n=4`$, $`(\mathrm{\Omega }_{2n}^0,J_\omega ,g_M)`$ is a Kähler manifold (but $`\theta `$ does not need to be constant). For any $`n1`$ and $`\theta `$ constant, $`F^{}\omega `$ is parallel, i.e., $`(M,J_\omega ,g_M)`$ is a Kähler manifold.
## 4 Minimal immersions with equal Kähler angles
Let us assume that $`F:MN`$ is minimal with equal Kähler angles. On a open set of $`M`$ where a orthonormal frame $`\{X_\alpha ,Y_\alpha =J_\omega (X_\alpha )\}`$ is defined, we have from (3.11) and (2.4), for any $`p`$, $`ZT_pM`$ and $`\mu ,\gamma `$,
$$2\mathrm{cos}\theta _Z\mu ,\gamma =i_ZF^{}\omega (\mu ,\gamma )=ig_Z\mu \gamma ig_Z\gamma \mu .$$
(4.1)
Note that $`F^{}\omega (_Z\mu ,\overline{\gamma })`$ $`=i\mathrm{cos}\theta _Z\mu ,\overline{\gamma }`$ $`=i\mathrm{cos}\theta \mu ,_Z\overline{\gamma }`$ $`=F^{}\omega (\mu ,_Z\overline{\gamma })`$. Hence, if $`\mu \gamma `$, $`_ZF^{}\omega (\mu ,\overline{\gamma })=d(F^{}\omega (\mu ,\overline{\gamma }))(Z)=0`$. Thus
$$g_Z\mu \overline{\gamma }=g_Z\overline{\gamma }\mu ,\mu \gamma $$
(4.2)
From (3.8), for each $`\mu `$,
$$\frac{i}{2}d\mathrm{cos}\theta (Z)=_ZF^{}\omega (\mu ,\overline{\mu })=g_Z\mu \overline{\mu }g_Z\overline{\mu }\mu (\text{no sumation over }\mu )$$
(4.3)
From (1.6), on $`M(𝒞)`$
$`\mathrm{}\kappa `$ $`=`$ $`4i{\displaystyle \underset{\beta }{}}Ricci^N(JdF(\beta ),dF(\overline{\beta }))`$ (4.6)
$`+{\displaystyle \frac{32}{\mathrm{sin}^2\theta }}{\displaystyle \underset{\beta ,\mu }{}}Im\left(R^N(dF(\beta ),dF(\mu ),dF(\overline{\beta }),JdF(\overline{\mu })+i\mathrm{cos}\theta dF(\overline{\mu }))\right)`$
$`{\displaystyle \frac{128\mathrm{cos}\theta }{\mathrm{sin}^4\theta }}{\displaystyle \underset{\beta ,\mu ,\rho }{}}Re\left(g_\beta \mu \overline{\rho }g_{\overline{\beta }}\rho \overline{\mu }\right)`$
$`+{\displaystyle \frac{64\mathrm{cos}\theta }{\mathrm{sin}^2\theta }}{\displaystyle \underset{\beta ,\mu ,\rho }{}}\left(|_\beta \mu ,\rho |^2+|_{\overline{\beta }}\mu ,\rho |^2\right),`$
where now $`\kappa =n\mathrm{log}\left(\frac{1+\mathrm{cos}\theta }{1\mathrm{cos}\theta }\right)`$. Since $`R(X,Y,Z,JW)`$ is skew-symmetric on $`(X,Y)`$ and symmetric on $`(Z,W)`$, $`_{\mu ,\beta }R^N(dF(\beta ),dF(\mu ),dF(\overline{\beta }),JdF(\overline{\mu }))=0`$. Then, from the Gauss equation and minimality of $`F`$,
$`(4.4)`$ $`=`$ $`{\displaystyle \underset{\beta ,\mu }{}}{\displaystyle \frac{32}{\mathrm{sin}^2\theta }}Im\left(i\mathrm{cos}\theta R^N(dF(\beta ),dF(\mu ),dF(\overline{\beta }),dF(\overline{\mu }))\right)`$
$`=`$ $`{\displaystyle \frac{32\mathrm{cos}\theta }{\mathrm{sin}^2\theta }}{\displaystyle \underset{\beta ,\mu }{}}R^M(\beta ,\mu ,\overline{\beta },\overline{\mu })+g(_{}dF(\beta ,\overline{\mu }),_{}dF(\mu ,\overline{\beta })).`$
Using the unitary basis $`\{\frac{\sqrt{2}}{\mathrm{sin}\theta }\mathrm{\Phi }(\rho ),\frac{\sqrt{2}}{\mathrm{sin}\theta }\mathrm{\Phi }(\overline{\rho })\}`$ of the normal bundle,
$`{\displaystyle \frac{32\mathrm{cos}\theta }{\mathrm{sin}^2\theta }}{\displaystyle \underset{\beta ,\mu }{}}g(_{}dF(\beta ,\overline{\mu }),_{}dF(\mu ,\overline{\beta }))={\displaystyle \frac{64\mathrm{cos}\theta }{\mathrm{sin}^4\theta }}{\displaystyle \underset{\beta ,\mu ,\rho }{}}(|g_\beta \overline{\mu }\rho |^2+|g_\beta \overline{\mu }\overline{\rho }|^2)=`$ (4.7)
$`=`$ $`{\displaystyle \frac{64\mathrm{cos}\theta }{\mathrm{sin}^4\theta }}{\displaystyle \underset{\beta ,\mu ,\rho }{}}(|g_\beta \overline{\rho }\mu |^2+|g_{\overline{\mu }}\beta \overline{\rho }|^2)={\displaystyle \frac{128\mathrm{cos}\theta }{\mathrm{sin}^4\theta }}{\displaystyle \underset{\beta ,\mu ,\rho }{}}|g_\beta \overline{\rho }\mu |^2.`$
From (4.2) and (4.3),
$`{\displaystyle \underset{\beta ,\mu ,\rho }{}}Re\left(g_\beta \mu \overline{\rho }g_{\overline{\beta }}\rho \overline{\mu }\right)={\displaystyle \underset{\beta ,\mu }{}}{\displaystyle \underset{\rho \mu }{}}|g_\beta \overline{\rho }\mu |^2+{\displaystyle \underset{\beta ,\mu }{}}Re\left(g_\beta \mu \overline{\mu }g_{\overline{\beta }}\mu \overline{\mu }\right)`$
$`={\displaystyle \underset{\beta ,\mu ,\rho }{}}|g_\beta \overline{\rho }\mu |^2{\displaystyle \underset{\beta ,\mu }{}}|g_\beta \overline{\mu }\mu |^2+{\displaystyle \underset{\beta ,\mu }{}}Re\left(g_\beta \mu \overline{\mu }g_{\overline{\beta }}\mu \overline{\mu }\right)`$
$`={\displaystyle \underset{\beta ,\mu ,\rho }{}}|g_\beta \overline{\rho }\mu |^2{\displaystyle \underset{\beta ,\mu }{}}Re\left({\displaystyle \frac{i}{2}}d\mathrm{cos}\theta (\beta )g_{\overline{\beta }}\mu \overline{\mu }\right),`$
so
$$(4.7)+(4.5)=\frac{128\mathrm{cos}\theta }{\mathrm{sin}^4\theta }\underset{\beta ,\mu }{}Re\left(\frac{i}{2}d\mathrm{cos}\theta (\beta )g_{\overline{\beta }}\mu \overline{\mu }\right).$$
On the other hand, Proposition 3.2 and minimality of $`F`$ gives,
$$\underset{\beta ,\mu }{}\frac{4\mathrm{cos}\theta }{\mathrm{sin}^2\theta }Re\left(ig_\beta \overline{\mu }\mu \overline{\beta }\right)=\frac{1n}{4}\mathrm{log}\mathrm{sin}^2\theta =\frac{(n1)\mathrm{cos}\theta }{2\mathrm{sin}^2\theta }\mathrm{cos}\theta .$$
Consequentely,
$`{\displaystyle \frac{128\mathrm{cos}\theta }{\mathrm{sin}^4\theta }}{\displaystyle \underset{\beta ,\mu }{}}Re\left({\displaystyle \frac{i}{2}}d\mathrm{cos}\theta (\beta )g_{\overline{\beta }}\mu \overline{\mu }\right)={\displaystyle \frac{128\mathrm{cos}\theta }{\mathrm{sin}^4\theta }}{\displaystyle \underset{\beta ,\mu }{}}Re\left({\displaystyle \frac{i}{2}}d\mathrm{cos}\theta (\overline{\beta })g_\beta \overline{\mu }\mu \right)`$
$`={\displaystyle \frac{64\mathrm{cos}\theta }{\mathrm{sin}^4\theta }}d\mathrm{cos}\theta (Re({\displaystyle \underset{\beta ,\mu }{}}ig_\beta \overline{\mu }\mu \overline{\beta }\left)\right)={\displaystyle \frac{8(n1)\mathrm{cos}\theta }{\mathrm{sin}^4\theta }}\mathrm{cos}\theta ^2.`$
That is,
$$(4.7)+(4.5)=\frac{8(n1)\mathrm{cos}\theta }{\mathrm{sin}^4\theta }\mathrm{cos}\theta ^2.$$
(4.8)
Using (3.9),
$`_{}J_\omega ^2`$ $`=`$ $`{\displaystyle \underset{\beta }{}}4_\beta J_\omega ,_{\overline{\beta }}J_\omega ={\displaystyle \underset{\beta }{}}{\displaystyle \underset{\mu ,\rho }{}}16\left(|_\beta J_\omega (\mu ),\rho |^2+|_\beta J_\omega (\overline{\mu }),\overline{\rho }|^2\right)`$ (4.9)
$`=`$ $`64{\displaystyle \underset{\beta ,\mu ,\rho }{}}\left(|_\beta \mu ,\rho |^2+|_{\overline{\beta }}\mu ,\rho |^2\right).`$
Thus we see that $`(4.6)=\frac{\mathrm{cos}\theta }{\mathrm{sin}^2\theta }_{}J_\omega ^2`$. So we have obtained the following formula:
###### Proposition 4.1
If $`N`$ is Kähler-Einstein with Ricci tensor $`Ricci^N=Rg`$, and $`F`$ is a minimal immersion with equal Kähler angles, on an open set without complex and Lagrangian points,
$`\mathrm{}\kappa `$ $`=`$ $`\mathrm{cos}\theta (2nR+{\displaystyle \frac{32}{\mathrm{sin}^2\theta }}{\displaystyle \underset{\beta ,\mu }{}}R^M(\beta ,\mu ,\overline{\beta },\overline{\mu })`$ (4.10)
$`+{\displaystyle \frac{1}{\mathrm{sin}^2\theta }}_{}J_\omega ^2+{\displaystyle \frac{8(n1)}{\mathrm{sin}^4\theta }}\mathrm{cos}\theta ^2).`$
Note that if $`n=1`$ we get the expression of Wolfson \[W\], $`\mathrm{}\kappa =2R\mathrm{cos}\theta `$.
###### Proposition 4.2
If $`N`$ is Kähler-Einstein with Ricci tensor $`Ricci^N=Rg`$, and $`F`$ is a minimal imersion with equal Kähler angles, then:
(i) If $`n=2`$,
$$_MnR\mathrm{sin}^2\theta \mathrm{cos}^2\theta Vol_M=0.$$
(4.11)
(ii) If $`n3`$ and $`F`$ has no complex points,
$$_MnR\mathrm{sin}^2\theta \mathrm{cos}^2\theta Vol_M=_M(n2)(n2+2\mathrm{cot}^2\theta )\mathrm{cos}\theta ^2Vol_M.$$
(4.12)
Proof. Multiplying (4.10) by $`\mathrm{sin}^2\theta \mathrm{cos}\theta `$, we get, on $`M𝒞`$, and using Lemma 3.2,
$`\mathrm{sin}^2\theta \mathrm{cos}\theta \mathrm{}\kappa `$ $`=`$ $`2n\mathrm{sin}^2\theta \mathrm{cos}^2\theta R+2SF^{}\omega ,F^{}\omega `$
$`+\mathrm{cos}^2\theta _{}J_\omega ^2+{\displaystyle \frac{8(n1)\mathrm{cos}^2\theta }{\mathrm{sin}^2\theta }}\mathrm{cos}\theta ^2.`$
On the other hand, $`\kappa =n\mathrm{log}\left(\frac{1+\mathrm{cos}\theta }{1\mathrm{cos}\theta }\right)`$, and so, $`\mathrm{}\kappa =\frac{2n}{\mathrm{sin}^2\theta }\mathrm{}\mathrm{cos}\theta +\frac{4n\mathrm{cos}\theta }{\mathrm{sin}^4\theta }\mathrm{cos}\theta ^2`$. Hence,
$`2n\mathrm{cos}\theta \mathrm{}\mathrm{cos}\theta +{\displaystyle \frac{4n\mathrm{cos}^2\theta }{\mathrm{sin}^2\theta }}\mathrm{cos}\theta ^2=`$
$`=`$ $`2n\mathrm{sin}^2\theta \mathrm{cos}^2\theta R+2SF^{}\omega ,F^{}\omega +\mathrm{cos}^2\theta _{}J_\omega ^2+{\displaystyle \frac{8(n1)\mathrm{cos}^2\theta }{\mathrm{sin}^2\theta }}\mathrm{cos}\theta ^2.`$
Recall that, from (3.1), and considering $`F^{}\omega `$ a 2-form, $`_{}F^{}\omega ^2=\frac{1}{2}\mathrm{cos}^2\theta _{}J_\omega ^2+n\mathrm{cos}\theta ^2`$. Since $`\mathrm{}\mathrm{cos}^2\theta =2\mathrm{cos}\theta \mathrm{}\mathrm{cos}\theta +2\mathrm{cos}\theta ^2`$, substituting this into (4.13), we have
$$n\mathrm{}\mathrm{cos}^2\theta =2n\mathrm{sin}^2\theta \mathrm{cos}^2\theta R+2SF^{}\omega ,F^{}\omega +2_{}F^{}\omega ^2+\frac{4(n2)\mathrm{cos}^2\theta }{\mathrm{sin}^2\theta }\mathrm{cos}\theta ^2$$
(4.14)
and, for $`n=2`$,
$`n\mathrm{}\mathrm{cos}^2\theta `$ $`=`$ $`2n\mathrm{sin}^2\theta \mathrm{cos}^2\theta R+2SF^{}\omega ,F^{}\omega +2_{}F^{}\omega ^2.`$ (4.15)
Let us now suppose that $`n3`$. Then, under the condition of no complex points, (4.14) is valied on $`\mathrm{\Omega }_{2n}^0`$ and also on $`\mathrm{\Omega }_0^0`$. From smoothness over all $`M`$ of all maps into consideration (the first three terms of the right-hand side of (4.14) are smooth, and the last term is also smooth for $`n2`$), and the fact that the set $`M(\mathrm{\Omega }_0^0\mathrm{\Omega }_0^{2n})`$ is a set of Lagrangian points with no interior, formula (4.14) is valid on all $`M`$. Integrating over $`M`$, and using (3.17), we have
$$_M2nR\mathrm{sin}^2\theta \mathrm{cos}^2\theta Vol_M=_M\left(2(n(n2)^2)+\frac{4(n2)\mathrm{cos}^2\theta }{\mathrm{sin}^2\theta }+2n\right)\mathrm{cos}\theta ^2Vol_M,$$
leading to (4.12). If $`n=2`$, we see that (4.15) is also valid at Lagrangian and complex points. In fact (see Lemma 3.2 and (3.1)), all terms of (4.15) vanish at interior points of the Lagrangian and complex sets. Since they are smooth on all $`M`$, they must vanish at boundary points of its complementary in $`M`$. Thus, the above equation is valid on all $`M`$, with or without complex or Lagrangian points, and all its terms are smooth. Then, (4.11) follows by integration on $`M`$ of (4.15), and use of (3.17). $`\mathbf{}`$
Proof of Theorem 1.2. and Theorem 1.3 If $`n=2`$ and $`R0`$, $`(4.11)`$ implies $`\mathrm{sin}^2\theta \mathrm{cos}^2\theta =0`$. Hence $`F`$ is either Lagrangian or a complex submanifold. If $`n3`$, and $`F`$ has no complex points, the right-hand side of (4.12) is non-negative, while the left-hand side is non-positive for $`R<0`$. Then, $`\mathrm{sin}^2\theta \mathrm{cos}^2\theta =0`$ must hold on all $`M`$, that is, $`F`$ is Lagrangian. If $`R=0`$, the right-hand side of (4.12) must vanish. Then, for $`n3`$, $`\mathrm{cos}\theta `$ must be constant, and we have proved Theorem 1.2. If $`\mathrm{cos}\theta `$ is constant, and if $`F`$ is not a complex submanifold, the right-hand side of (4.12) vanishes. Hence, if $`R0`$, $`F`$ is Lagrangian, and Theorem 1.3 is proved. $`\mathbf{}`$
Proof of Theorem 1.4. If $`M`$ is not Lagrangian, under the curvature condition on $`M`$, by Proposition 3.4, for $`n=2`$, or $`3`$, $`(M,J_\omega ,g_M)`$ is a Kähler manifold and $`\mathrm{cos}\theta `$ is constant. So, if $`M`$ is not a complex submanifold, it has no complex directions, and by (4.11), or (4.12), $`R=0`$. In general, if $`n1`$ and $`\theta `$ is constant, Proposition 3.4 also applies. $`\mathbf{}`$
Under the conditions of Theorem 1.4, if $`M`$ is homeomorphic to a 4 or a 6 dimensional sphere, minimaly immersed into a Kähler-Einstein manifold, and with equal Kähler angles, then it must be Lagrangian, for it is well known that such manifolds cannot carry a Kähler structure. Obviously, any Riemannian manifold $`M`$ with strictly positive isotropic scalar curvature cannot carry any Kähler structure. Moreover, such condition for $`n=2`$ would imply $`M`$ to be homeomorphic to a 4-sphere. We also remark that we only need to require $`S_{isot}(\{Z_\alpha \})0`$ on the maximal totally isotropic subspace $`\{Z_\alpha \}`$ defined by a diagonalizing orthonormal basis of $`F^{}\omega `$, and outside Lagrangian points, to obtain the same conclusion given in Theorem 1.4.
As an observation, Theorem 1.4 should be compared with the following lemma:
###### Lemma 4.1
Let $`F`$ be a minimal immersion, and $`n2`$. If $`\mathrm{cos}\theta `$ is constant $`1,0`$, then
$`(i)`$ $`(A,B,C)g_ABC`$ is symmetric whenever $`A,B`$, and $`C`$ are not all of the same type.
$`(ii)`$ $`_{\overline{\beta }}\mu ,\gamma =0,\beta ,\mu ,\gamma `$.
$`(iii)`$ $`F^{}\omega `$ is an harmonic 2-form.
$`(iv)`$ $`32_{\beta ,\mu }R^M(\beta ,\mu ,\overline{\beta },\overline{\mu })=64_{\beta ,\mu ,\rho }|_\beta \mu ,\rho |^2=_{}J_\omega ^20.`$
Proof. Since $`\mathrm{cos}\theta `$ is constant, we obtain $`(4.3)=0`$. This, together (4.2), and the symmetry of $`_{}dF`$, proves $`(i)`$. But $`(i)`$ and (4.1) implys $`(ii)`$. $`(iii)`$ comes from (3.16). Now we prove $`(vi)`$. Since $`F^{}\omega `$ is harmonic, from Weitzenböck formula (3.14) we conclude $`SF^{}\omega ,F^{}\omega =_{}F^{}\omega ^2`$. Lemma 3.2 and (3.1) (but considering $`F^{}\omega `$ a 2-form) gives $`(iv)`$. $`\mathbf{}`$
Remark 3. If $`N`$ is a Kähler manifold of constant holomorphic sectional curvature equal to $`K`$ ( and so $`R=\frac{(2n+1)K}{2}`$), and the isotropic scalar curvature of $`M`$ satisfies $`S_{isot}c=constant`$, we get from Gauss equation, with $`\{X_\alpha ,Y_\alpha \}`$ a diagonalizing orthonormal basis of $`F^{}\omega `$,
$$\underset{\rho ,\mu }{}R^M(\mu ,\rho ,\overline{\mu },\overline{\rho })=\frac{n(n1)}{16}\mathrm{sin}^2\theta K\underset{\rho ,\mu }{}_{}dF(\mu ,\overline{\rho })^2,$$
(4.16)
that $`c\frac{n(n1)K}{4}`$. Thus, non-negative isotropic scalar curvature on $`M`$ is a possible condition for $`K0`$. In the case $`K=0`$, that is, $`N`$ is the flat complex torus, then (4.16) (with $`K=0`$) is valied for any orthonormal basis $`\{X_\alpha ,Y_\alpha \}`$. This implies that, for $`n2`$, $`F`$ must be totally geodesic, and so $`M`$ is flat.
We also note that if $`c=\frac{nR}{4}`$, the right-hand side of (4.10) becomes $`>0`$, outside Lagrangian points. An application of the maximum principle at a maximum point of $`\kappa `$ would conclude that $`F`$ must be Lagrangian. But such a lower bound $`c`$ is not possible for the scalar isotropic curvature of $`M`$ minimaly immersed in $`N`$ with constant holomorphic sectional curvature $`K>0`$.
Remark 3. If $`n2`$ and $`F`$ is a pluriminimal immersion with equal Kähler angles into a Kähler-Einstein manifold $`N`$, and $`F`$ is not a complex submanifold, then $`N`$ must be Ricci-flat. Moreover, since $`F`$ has constant equal Kähler angles, the scalar isotropic curvature of $`M`$ with respect to the maximal isotropic subspace defined by a diagonalizing orthonormal basis of $`F^{}\omega `$ will be $`0`$, with equality to zero iff $`(M,J_\omega ,g_M)`$ is Kähler (see Lemma 4.1). We leave the following question: Is $`(M,J_\omega ,g_M)`$ Kähler manifold a sufficient condition for a minimal immersion $`F`$, with constant equal Kähler angle, immersed into a Ricci-flat Kähler manifold $`N`$, to be pluriminimal? If $`N`$ is the flat complex torus and $`F:MN`$ is minimal, under the conditions stated in the question, the Gauss equation implies that $`F`$ is pluriminimal. A way to find pluriminimal submanifolds in hyper-Kähler manifolds is given in the next example, where the assumption of non-negative isotropic curvature does not imply necessarely $`F`$ totally geodesic (and $`M`$ flat), since hyper-Kähler manifolds do not need to be flat.
Example. Let $`(N,I,J,g)`$ be an hyper-Kähler manifold of real dimension $`8`$. Thus, $`I`$ and $`J`$ are two $`g`$-orthogonal complex structures on $`N`$, such that $`IJ=JI`$ and $`_{}I=_{}J=0`$, where $`_{}`$ is the Levi-Civita connection relative to $`g`$. It is known that such manifolds are Ricci-flat (\[B\]). Set $`K=IJ`$. For each $`\nu `$, $`\varphi `$, we take $`\mathrm{`}\mathrm{`}\nu \varphi \mathrm{"}=(\mathrm{cos}\nu ,\mathrm{sin}\nu \mathrm{cos}\varphi ,\mathrm{sin}\nu \mathrm{sin}\varphi )S^2`$, and define $`J_{\nu \varphi }=\mathrm{cos}\nu I+\mathrm{sin}\nu \mathrm{cos}\varphi J+\mathrm{sin}\nu \mathrm{sin}\varphi K`$. These $`J_{\nu \varphi }`$ are the complex structures on $`N`$ compatible with its hyper-Kähler structure, that is, they are $`g`$-orthogonal and $`_{}J_{\nu \varphi }=0`$.
Two of such complex structures, $`J_{\nu \varphi }`$ and $`J_{\mu \rho }`$, anti-commute at a point $`p`$ iff $`J_{\nu \varphi }(X)`$ and $`J_{\mu \rho }(X)`$ are orthogonal for some non-zero $`XT_pN`$, iff $`\nu \varphi `$ and $`\mu \rho `$ are orthogonal in $`IR^3`$. Thus, they anti-commute at a point $`p`$ iff they anti-commute everywhere. If that is the case $`J_{\nu \varphi }J_{\mu \rho }=J_{\sigma ϵ}`$, where $`\{\nu \varphi ,\mu \rho ,\sigma ϵ\}`$ is a direct orthonormal basis of $`IR^3`$. For each unit vector $`XT_pN`$, set $`H_X=span\{X,IX,JX,KX\}`$ $`=span\{X,J_{\nu \varphi }(X),J_{\mu \rho }(X),J_{\sigma ϵ}(X)\}`$, for any orthonormal basis $`\{\nu \varphi ,\mu \rho ,\sigma ϵ\}`$. If $`YH_X^{}`$ is another unit vector, then $`H_XH_Y`$. Let $`\omega _{\nu \varphi }`$ be the Kähler form of $`(N,J_{\nu \varphi },g)`$. Let $`E`$ be a 4-dimensional vector sub-space of $`T_pN`$. We first note that $`E=H_X`$ for some $`XE`$, iff $`J_{\nu \varphi }(E)E`$ for any $`\nu ,\varphi `$. If that is the case, then $`E`$ is not a Lagrangian subspace with respect to any complex structure $`J_{\mu \rho }`$. In general, $`E`$ contains a $`J_{\nu \varphi }`$-complex line for some $`\nu \varphi `$ iff $`dim(EH_X)2`$ for some $`XE`$. If that is the case, and if $`E`$ is a Lagrangian subspace of $`T_pN`$ with respect to $`J_{\mu \rho }`$, then $`\nu \varphi \mu \rho `$. Furthermore, if $`E`$ is a $`J_{\nu \varphi }`$-complex subspace, then $`E`$ is $`J_{\mu \rho }`$-Lagrangian iff there exist an orthonormal basis $`\{X,J_{\nu \varphi }X,Y,J_{\nu \varphi }Y\}`$ of $`E`$ with $`H_XH_Y`$. To see this, let us suppose $`E`$ is $`J_{\nu \varphi }`$-complex subspace and $`J_{\mu \rho }`$-Lagrangian. We take $`\{X,J_{\nu \varphi }X,Y,J_{\nu \varphi }Y\}`$ an ortonormal basis of $`E`$. Then $`Yspan\{X,J_{\nu \varphi }X,J_{\mu \rho }X\}^{}`$. So $`Y=tJ_{\sigma ϵ}X+\stackrel{~}{Y}`$, for some $`tIR`$ and $`\stackrel{~}{Y}H_X^{}`$, and where $`\{\nu \varphi ,\mu \rho ,\sigma ϵ\}`$ is an ortonormal basis of $`IR^3`$. As $`EH_X`$, $`\stackrel{~}{Y}0`$. From $`0=J_{\mu \rho }Y,J_{\nu \varphi }X`$, we get $`t=0`$. Thus, $`YH_X^{}`$. We observe that, in general, $`J_{\mu \rho }`$-Lagrangian subspaces do not need to be $`J_{\nu \varphi }`$-complex, as for example $`E=\{X,J_{\nu \varphi }X,Y,J_{\sigma ϵ}Y\}`$, with $`YH_X^{}`$, that contains two orthogonal complex lines for different complex strutures.
Any $`J_{\nu \varphi }`$-complex submanifold $`F:MN`$ of real dimension $`4`$, such that, for each point $`pM`$, there exist an orthonormal basis $`\{X,J_{\nu \varphi }X,Y,J_{\nu \varphi }Y\}`$ of $`T_pM`$ with $`H_XH_Y`$, is, for each $`\mu \rho `$, a minimal submanifold of $`(N,J_{\mu \rho },g)`$ with constant equal Kähler angles, and $`\pm J_{\nu \varphi }`$ is also the complex structure of $`M`$ which comes from polar decomposition of $`\omega _{\mu \rho }`$ restricted to $`M`$. In fact, such an orthonormal basis of $`T_pM`$ diagonalizes $`\omega _{\mu \rho }`$ restricted to $`M`$, and the Kähler angle $`\theta `$ is such that $`\mathrm{cos}\theta =\pm \nu \varphi ,\mu \rho `$, where $`<,>`$ is the inner product of $`IR^3`$. Next proposition is an application of Theorem 1.4, for 4-dimensional submanifolds of $`N`$, where $`\omega _I`$ is the Kähler form of $`(N,I,g)`$:
###### Proposition 4.3
Let $`F:MN`$ be a minimal immersion of a compact, oriented 4-dimensional submanifold with non-negative isotropic scalar curvature, and such that $`\nu \varphi S^2`$, $`F`$ has equal Kähler angles with respect to $`J_{\nu \varphi }`$. If $`pM`$ and $`XT_pM`$, unit vector, such that $`dim(T_pMH_X)2`$, then there exists $`\nu \varphi S^2`$ such that $`M`$ is a $`J_{\nu \varphi }`$-complex submanifold. Furthermore, if $`J_{\nu \varphi }=I`$ then $`F:M(N,I,g)`$ is obviously pluriminimal. If $`J_{\nu \varphi }I`$ but $`T_pMH_X^{}\{0\}`$, then $`F^{}\omega _I=\mathrm{cos}\nu J_{\nu \varphi }`$, and if $`F`$ is not $`J_I`$-Lagrangian, $`F:M(N,I,g)`$ is still pluriminimal.
Note that, if $`T_pM=H_X`$, then $`J_{\nu \varphi }`$ can be chosen equal to $`I`$. The first conclusion of this result is the 4-dimensional version of a result of Wolfson \[W\], for $`M`$ a real surface and $`N`$ a Ricci-flat K3 surface. In the latter case, there is only one Kähler angle, $`X`$ $`dim(T_pMH_X)=2`$ is automatically satisfied, and the isotropic scalar curvature is always zero.
Proof. From the assumption, $`dim(T_pMH_X)2`$, we may take a unit vector $`ZT_pMH_X`$ such that $`ZX`$. Then, $`Z=J_{\nu \varphi }(X)`$ for some $`\nu \varphi `$. Thus, $`span\{X,J_{\nu \varphi }(X)\}T_pM`$. This implies $`F^{}\omega _{\nu \varphi }(X,J_{\nu \varphi }(X))=1`$. As the Kähler angles are equal, $`\mathrm{cos}\theta _{\nu \varphi }=1`$ at $`p`$. Applying Theorem 1.4 to $`F:M(N,J_{\nu \varphi },g)`$, $`F^{}\omega _{\nu \varphi }=\mathrm{cos}\theta _{\nu \varphi }J_{\omega _{\nu \varphi }}`$ with $`\mathrm{cos}\theta _{\nu \varphi }`$ constant. Then $`\mathrm{cos}\theta _{\nu \varphi }=1`$ everywhere. That is, $`M`$ is a $`J_{\nu \varphi }`$-complex submanifold. Moreover, from the second assumption, $`T_pMH_X^{}\{0\}`$, we may take a unit vector $`YT_pMH_X^{}`$. Then $`\{X,J_{\nu \varphi }X,Y,J_{\nu \varphi }Y\}`$ constitutes an orthonormal basis of $`T_pM`$, that diagonalizes $`F^{}\omega _I`$, and $`F^{}\omega _I=\mathrm{cos}\nu J_{\nu \varphi }`$. This means that $`\nu `$ or $`\nu +\pi `$ is the constant Kähler angle of $`F:M(N,I,g)`$, and, since $`M`$ is a $`J_{\nu \varphi }`$-complex submanifold, it is pluriharmonic with respect to $`\pm J_{\nu \varphi }`$, and so, if $`\mathrm{cos}\nu 0`$, it is pluriminimal as an immersion into $`(N,I,g)`$. $`\mathbf{}`$
## 5 Appendix: The computation of $`\mathrm{}\kappa `$
We prove (1.6) for $`F`$ minimal and outside complex and Lagrangian points. First, we compute some derivative formulas of a determinant, which we will need.
###### Lemma 5.1
Let $`A:M_{m\times m}(lC)`$ be a smooth map of matrices $`pA(p)=[A_1,\mathrm{},A_m]`$, where $`A_i(p)`$ is a column vector of $`lC^m`$ and $`M`$ is a Riemannian manifold with its Levi-Civita connection $`_{}`$. Assume that, at $`p_0`$, $`A(p_0)`$ is a diagonal matrix $`D=D(\lambda _1,\mathrm{},\lambda _m)`$. Then, at $`p_0`$
$$d(\text{det }A)(Z)=\underset{1jm}{}\left(\underset{kj}{}\lambda _k\right)dA_j^j(Z),$$
$`\text{Hess }(detA)(Z,W)=_{}d(detA)(Z,W)=`$
$`=`$ $`{\displaystyle \underset{1j,km}{}}({\displaystyle \underset{sj,k}{}}\lambda _s)\text{det}\left[\begin{array}{cc}dA_j^j(Z)& dA_j^k(Z)\\ dA_k^j(W)& dA_k^k(W)\end{array}\right]+{\displaystyle \underset{1jm}{}}({\displaystyle \underset{sj}{}}\lambda _s)\text{Hess }A_j^j(Z,W).`$
In particular, if $`e_1,\mathrm{},e_r`$ is an orthonormal basis of $`T_{p_0}M`$, then, at $`p_0`$,
$`\mathrm{}(detA)=Trace\text{Hess }(detA)=`$
$`=`$ $`{\displaystyle \underset{1\alpha r}{}}{\displaystyle \underset{1j,km}{}}({\displaystyle \underset{sj,k}{}}\lambda _s)\text{det}\left[\begin{array}{cc}dA_j^j(e_\alpha )& dA_j^k(e_\alpha )\\ dA_k^j(e_\alpha )& dA_k^k(e_\alpha )\end{array}\right]+{\displaystyle \underset{1jm}{}}\left({\displaystyle \underset{sj}{}}\lambda _s\right)\mathrm{}A_j^j.`$
On each $`\mathrm{\Omega }_{2k}^0`$, the complex structure $`J_\omega `$ and the sub-vector bundle $`𝒦_\omega ^{}`$ are smooth. Moreover, $`J_\omega `$ is $`g_M`$-orthogonal. Thus, for each $`p_0\mathrm{\Omega }_{2k}^0`$, there exists a locally $`g_M`$-orthonormal frame of $`𝒦_\omega ^{}`$ defined on a neighbourhood of $`p_0`$, of the form $`X_1,J_\omega X_1,\mathrm{},X_k,J_\omega X_k`$. We enlarge this frame to a $`g_M`$-orthonormal local frame on $`M`$, on a neighbourhood of $`p_0`$:
$$X_1,Y_1=J_\omega X_1,\mathrm{},X_k,Y_k=J_\omega X_k,X_{k+1},Y_{k+1},\mathrm{},X_n,Y_n$$
(5.3)
where $`X_{k+1},Y_{k+1},\mathrm{}X_n,Y_n`$ is any $`g_M`$-orthonormal frame of $`𝒦_\omega `$, and which at $`p_0`$ is a diagonalizing basis of $`F^{}\omega `$. Note that in general it is not possible to get smooth diagonalizing $`g_M`$-orthonormal frames in a whole neighbourhood of a point $`p_0`$, unless , for instance, $`F^{}\omega `$ has equal Kähler angles. We use the notations in section 3.1. We define a local complex structure on a neighbourhood of $`p_0\mathrm{\Omega }_{2k}^0`$ as $`\stackrel{~}{J}=J_\omega J^{}`$, where $`J_\omega `$ is defined only on $`𝒦_\omega ^{}`$, and $`J^{}`$ is the local complex structure on $`𝒦_\omega `$, defined on a neighbourhood of $`p_0`$ by
$$J^{}Z_\alpha =iZ_\alpha ,J^{}Z_{\overline{\alpha }}=iZ_{\overline{\alpha }},\alpha k+1.$$
(5.4)
Thus, the vectors $`Z_\alpha `$ are of type (1,0) with respect to $`\stackrel{~}{J}`$, for $`\alpha `$. Since $`\stackrel{~}{J}`$ is $`g_M`$-orthogonal, then, $`\alpha ,\beta `$, on a neighbourhood of $`p_0`$,
$$_Z\stackrel{~}{J}(\alpha ),\beta =2i_Z\alpha ,\beta =\alpha ,_Z\stackrel{~}{J}(\beta ),_Z\stackrel{~}{J}(\alpha ),\overline{\beta }=0,$$
(5.5)
Note that $`F^{}\omega `$ and $`\stackrel{~}{g}`$, where $`\stackrel{~}{g}`$ is given in (1.1), as 2-tensors, are both of type $`(1,1)`$ with respect to $`\stackrel{~}{J}`$, and have the same kernel $`𝒦_\omega `$. They are related by $`\stackrel{~}{g}(X,Y)=F^{}\omega (X,J_\omega Y)=F^{}\omega (X,\stackrel{~}{J}Y)`$. Set $`\stackrel{~}{g}_{AB}=\stackrel{~}{g}(A,B)`$, and define $`\overline{\overline{B}}=B`$, $`A,B\{1,\mathrm{},n,\overline{1},\mathrm{},\overline{n}\}`$, and set $`ϵ_\alpha =+1`$, $`ϵ_{\overline{\alpha }}=1`$, $`1\alpha n`$. Let $`1\alpha ,\beta n`$, $`A,B\{1,\mathrm{},n,\overline{1},\mathrm{},\overline{n}\}`$, and $`C\{1,\mathrm{},n\}\{\overline{k+1},\mathrm{},\overline{n}\}`$. Then
$$\begin{array}{cc}F^{}\omega (\alpha ,C)=g(JdF(\alpha ),dF(C))=0\hfill & p\text{ near }p_0\hfill \\ F^{}\omega (\alpha ,\overline{\beta })=g(JdF(\alpha ),dF(\overline{\beta }))=\frac{i}{2}\delta _{\alpha \beta }\mathrm{cos}\theta _\alpha \hfill & \text{at }p_0\hfill \\ \stackrel{~}{g}_{AB}=iϵ_BF^{}\omega (A,B)=iϵ_Bg(JdF(A),dF(B))\hfill & p\text{ near }p_0\hfill \\ \stackrel{~}{g}_{\alpha C}=\stackrel{~}{g}_{\overline{\alpha }\overline{C}}=0\hfill & p\text{ near }p_0\hfill \\ \stackrel{~}{g}_{\alpha \overline{\beta }}=\stackrel{~}{g}_{\overline{\alpha }\beta }=\frac{1}{2}\delta _{\alpha \beta }\mathrm{cos}\theta _\alpha \hfill & \text{at }p_0\hfill \end{array}\}.$$
(5.6)
At a point $`p_0`$, with Kähler angles $`\theta _\alpha `$, $`g_M\pm \stackrel{~}{g}`$ is represented in the unitary basis $`\{\sqrt{2}\alpha ,\sqrt{2}\overline{\alpha }\}`$, by the diagonal matrix $`g_M\pm \stackrel{~}{g}=D(1\pm \mathrm{cos}\theta _1,\mathrm{},1\pm \mathrm{cos}\theta _n,1\pm \mathrm{cos}\theta _1,\mathrm{},1\pm \mathrm{cos}\theta _n)`$, and so
$$det(g_M\pm \stackrel{~}{g})=\underset{1\alpha n}{}(1\pm \mathrm{cos}\theta _\alpha )^2.$$
(5.7)
If $`p_0`$ is a point without complex directions, $`\mathrm{cos}\theta _\alpha 1`$, $`\alpha \{1,\mathrm{},n\}`$, then $`\stackrel{~}{g}<g_M`$. Thus, on a neighbourwood of $`p_0`$, we may consider the map $`\kappa `$.
$$\kappa =\frac{1}{2}\mathrm{log}\left(\frac{det(g_M+\stackrel{~}{g})}{det(g_M\stackrel{~}{g})}\right)=\underset{1\alpha n}{}\mathrm{log}\left(\frac{1+\mathrm{cos}\theta _\alpha }{1\mathrm{cos}\theta _\alpha }\right).$$
(5.8)
This map is continuous outside the complex points, and smooth on each $`\mathrm{\Omega }_{2k}^0`$. We wish to compute $`\mathrm{}\kappa `$ on $`\mathrm{\Omega }_{2k}^0`$.
###### Lemma 5.2
At $`p_0\mathrm{\Omega }_{2k}^0`$, without complex directions and for $`Z,WT_{p_0}M`$,
$$d(det(g_M\pm \stackrel{~}{g}))(Z)=\pm 4\underset{1\mu n}{}\frac{_{1\alpha n}(1\pm \mathrm{cos}\theta _\alpha )^2}{(1\pm \mathrm{cos}\theta _\mu )}d\stackrel{~}{g}_{\mu \overline{\mu }}(Z),$$
$`Hess\left(det(g_M\pm \stackrel{~}{g})\right)(Z,W)=`$
$`=`$ $`16\left({\displaystyle \underset{1\alpha n}{}}(1\pm \mathrm{cos}\theta _\alpha )^2\right){\displaystyle \underset{\mu ,\rho }{}}{\displaystyle \frac{1}{(1\pm \mathrm{cos}\theta _\mu )(1\pm \mathrm{cos}\theta _\rho )}}d\stackrel{~}{g}_{\mu \overline{\mu }}(Z)d\stackrel{~}{g}_{\rho \overline{\rho }}(W)`$
$`8\left({\displaystyle \underset{1\alpha n}{}}(1\pm \mathrm{cos}\theta _\alpha )^2\right){\displaystyle \underset{\mu ,\rho }{}}{\displaystyle \frac{1}{(1\pm \mathrm{cos}\theta _\mu )(1\pm \mathrm{cos}\theta _\rho )}}d\stackrel{~}{g}_{\mu \overline{\rho }}(W)d\stackrel{~}{g}_{\rho \overline{\mu }}(Z)`$
$`\pm 4\left({\displaystyle \underset{1\alpha n}{}}(1\pm \mathrm{cos}\theta _\alpha )^2\right){\displaystyle \underset{\mu }{}}{\displaystyle \frac{1}{(1\pm \mathrm{cos}\theta _\mu )}}Hess\stackrel{~}{g}_{\mu \overline{\mu }}(Z,W).`$
Proof. Using the unitary basis $`\{\sqrt{2}\alpha ,\sqrt{2}\overline{\alpha }\}`$ of $`T_p^cM`$, for $`p`$ near $`p_0`$, $`g_M+\stackrel{~}{g}`$ is represented by the matrix
$$g_M\pm \stackrel{~}{g}=\left[\begin{array}{cc}g_M\pm \stackrel{~}{g}(\sqrt{2}\alpha ,\sqrt{2}\overline{\gamma })& g_M\pm \stackrel{~}{g}(\sqrt{2}\alpha ,\sqrt{2}\gamma )\\ g_M\pm \stackrel{~}{g}(\sqrt{2}\overline{\alpha },\sqrt{2}\overline{\gamma })& g_M\pm \stackrel{~}{g}(\sqrt{2}\overline{\alpha },\sqrt{2}\gamma )\end{array}\right]=\left[\begin{array}{cc}\delta _{\alpha \gamma }\pm 2\stackrel{~}{g}_{\alpha \overline{\gamma }}& 0\\ 0& \delta _{\alpha \gamma }\pm 2\stackrel{~}{g}_{\overline{\alpha }\gamma }\end{array}\right]$$
that at $`p_0`$ is the diagonal matrix $`D(1\pm \mathrm{cos}\theta _1,\mathrm{},1\pm \mathrm{cos}\theta _n,1\pm \mathrm{cos}\theta _1,\mathrm{},1\pm \mathrm{cos}\theta _n)`$. The lemma follows as a simple application of lemma 5.1, and noting that $`\stackrel{~}{g}_{\mu \overline{\rho }}=\stackrel{~}{g}_{\overline{\rho }\mu }`$. $`\mathbf{}`$
On $`\mathrm{\Omega }_{2k}^0`$,
$`2\mathrm{}\kappa `$ $`=`$ $`\mathrm{}\mathrm{log}(det(g_M+\stackrel{~}{g}))\mathrm{}\mathrm{log}(det(g_M\stackrel{~}{g}))`$
$`=`$ $`{\displaystyle \frac{\mathrm{}(det(g_M+\stackrel{~}{g}))}{det(g_M+\stackrel{~}{g})}}{\displaystyle \frac{d(det(g_M+\stackrel{~}{g}))^2}{(det(g_M+\stackrel{~}{g}))^2}}{\displaystyle \frac{\mathrm{}(det(g_M\stackrel{~}{g}))}{det(g_M\stackrel{~}{g})}}+{\displaystyle \frac{d(det(g_M\stackrel{~}{g}))^2}{(det(g_M\stackrel{~}{g}))^2}}.`$
From the above lemma and
$`d(det(g_M\pm \stackrel{~}{g}))^2`$ $`=`$ $`4{\displaystyle \underset{\beta }{}}d(det(g_M\pm \stackrel{~}{g}))(\beta )d(det(g_M\pm \stackrel{~}{g}))(\overline{\beta })`$
$`\mathrm{}det(g_M\pm \stackrel{~}{g})`$ $`=`$ $`4{\displaystyle \underset{\beta }{}}Hess(det(g_M\pm \stackrel{~}{g}))(\beta ,\overline{\beta })`$
we have at $`p_0`$,
$$2\mathrm{}\kappa =\underset{\beta ,\mu ,\rho }{}\frac{64(\mathrm{cos}\theta _\mu +\mathrm{cos}\theta _\rho )}{\mathrm{sin}^2\theta _\mu \mathrm{sin}^2\theta _\rho }d\stackrel{~}{g}_{\mu \overline{\rho }}(\overline{\beta })d\stackrel{~}{g}_{\rho \overline{\mu }}(\beta )+\underset{\beta ,\mu }{}\frac{32}{\mathrm{sin}^2\theta _\mu }Hess\stackrel{~}{g}_{\mu \overline{\mu }}(\beta ,\overline{\beta }).$$
(5.9)
Recalling (2.4), and $`d(F^{}\omega (X,Y))(Z)=_ZF^{}\omega (X,Y)+F^{}\omega (_ZX,Y)+F^{}\omega (X,_ZY)`$, using (5.4), we obtain
###### Lemma 5.3
$`p`$ near $`p_0\mathrm{\Omega }_{2k}^0`$, $`ZT_p^cM`$, and $`\mu ,\gamma \{1,\mathrm{},n\}`$
$`d\stackrel{~}{g}_{\mu \overline{\gamma }}(Z)`$ $`=`$ $`ig_Z\mu \overline{\gamma }ig_Z\overline{\gamma }\mu +2{\displaystyle \underset{\rho }{}}\left(_Z\mu ,\overline{\rho }\stackrel{~}{g}_{\rho \overline{\gamma }}+_Z\overline{\gamma },\rho \stackrel{~}{g}_{\mu \overline{\rho }}\right)`$
$`0=d\stackrel{~}{g}_{\mu \gamma }(Z)`$ $`=`$ $`ig_Z\mu \gamma +ig_Z\gamma \mu +2{\displaystyle \underset{\rho }{}}\left(_Z\mu ,\rho \stackrel{~}{g}_{\overline{\rho }\gamma }_Z\gamma ,\rho \stackrel{~}{g}_{\mu \overline{\rho }}\right).`$
In particular, at $`p_0`$
$`d\stackrel{~}{g}_{\mu \overline{\gamma }}(Z)`$ $`=`$ $`ig_Z\mu \overline{\gamma }ig_Z\overline{\gamma }\mu (\mathrm{cos}\theta _\mu \mathrm{cos}\theta _\gamma )_Z\mu ,\overline{\gamma }`$
$`0=d\stackrel{~}{g}_{\mu \gamma }(Z)`$ $`=`$ $`ig_Z\mu \gamma +ig_Z\gamma \mu +(\mathrm{cos}\theta _\mu +\mathrm{cos}\theta _\gamma )_Z\mu ,\gamma .`$
###### Lemma 5.4
If F is minimal and $`p_0\mathrm{\Omega }_{2k}^0`$ is a point without complex directions, then for each $`\mu \{1,\mathrm{},n\}`$
$`{\displaystyle \underset{1\beta n}{}}Hess\stackrel{~}{g}_{\mu \overline{\mu }}(\beta ,\overline{\beta })={\displaystyle \underset{1\beta n}{}}d\left(d\stackrel{~}{g}_{\mu \overline{\mu }}(\beta )\right)(\overline{\beta })d\stackrel{~}{g}_{\mu \overline{\mu }}(_{\overline{\beta }}\beta )=`$
$`={\displaystyle \underset{1\beta n}{}}`$ $`iR^N(dF(\beta ),dF(\overline{\beta }),dF(\mu ),JdF(\overline{\mu })+i\mathrm{cos}\theta _\mu dF(\overline{\mu }))`$
$`+2Im\left(R^N(dF(\beta ),dF(\mu ),dF(\overline{\beta }),JdF(\overline{\mu })+i\mathrm{cos}\theta _\mu dF(\overline{\mu }))\right)`$
$`+2{\displaystyle \underset{1\rho n}{}}{\displaystyle \frac{(\mathrm{cos}\theta _\rho \mathrm{cos}\theta _\mu )}{\mathrm{sin}^2\theta _\rho }}\left(|g_\beta \mu \rho |^2+|g_\beta \overline{\mu }\overline{\rho }|^2\right)`$
$`2{\displaystyle \underset{1\rho n}{}}{\displaystyle \frac{(\mathrm{cos}\theta _\rho +\mathrm{cos}\theta _\mu )}{\mathrm{sin}^2\theta _\rho }}\left(|g_\beta \mu \overline{\rho }|^2+|g_\beta \overline{\mu }\rho |^2\right)`$
$`+{\displaystyle \underset{1\rho n}{}}2i_\mu \beta ,\overline{\rho }g_{\overline{\beta }}\rho \overline{\mu }2i_\mu \beta ,\rho g_{\overline{\beta }}\overline{\rho }\overline{\mu }2i_\mu \overline{\beta },\overline{\rho }g_\rho \beta \overline{\mu }`$
$`+{\displaystyle \underset{1\rho n}{}}2i_{\overline{\beta }}\mu ,\overline{\rho }g_\beta \rho \overline{\mu }2i_\mu \overline{\beta },\rho g_{\overline{\rho }}\beta \overline{\mu }+2i_{\overline{\beta }}\mu ,\rho g_{\overline{\rho }}\beta \overline{\mu }`$
$`+{\displaystyle \underset{1\rho n}{}}2i_{\overline{\mu }}\beta ,\overline{\rho }g_{\overline{\beta }}\rho \mu +2i_{\overline{\mu }}\beta ,\rho g_{\overline{\beta }}\overline{\rho }\mu +2i_{\overline{\mu }}\overline{\beta },\overline{\rho }g_\rho \beta \mu `$
$`+{\displaystyle \underset{1\rho n}{}}2i_{\overline{\beta }}\overline{\mu },\overline{\rho }g_\rho \beta \mu +2i_{\overline{\mu }}\overline{\beta },\rho g_{\overline{\rho }}\beta \mu 2i_{\overline{\beta }}\overline{\mu },\rho g_{\overline{\rho }}\beta \mu `$
$`+{\displaystyle \underset{1\rho n}{}}2i_{\overline{\beta }}\overline{\mu },\overline{\rho }g_\beta \mu \rho +2i_{\overline{\beta }}\overline{\mu },\rho g_\beta \mu \overline{\rho }2i_{\overline{\beta }}\mu ,\rho g_\beta \overline{\mu }\overline{\rho }`$
$`+{\displaystyle \underset{1\rho n}{}}2i_{\overline{\beta }}\mu ,\overline{\rho }g_\beta \overline{\mu }\rho +2i_\beta \mu ,\overline{\rho }g_{\overline{\beta }}\rho \overline{\mu }2i_\beta \mu ,\overline{\rho }g_{\overline{\beta }}\overline{\mu }\rho `$
$`+{\displaystyle \underset{1\rho n}{}}2i_\beta \overline{\mu },\rho g_{\overline{\beta }}\mu \overline{\rho }2i_\beta \overline{\mu },\rho g_{\overline{\beta }}\overline{\rho }\mu `$
$`2{\displaystyle \underset{1\rho n}{}}(\mathrm{cos}\theta _\mu \mathrm{cos}\theta _\rho )\left(|_\beta \mu ,\overline{\rho }|^2+|_{\overline{\beta }}\mu ,\overline{\rho }|^2\right).`$
Proof. We denote by $`_X_YdF`$ the covariant derivative of $`_YdF`$ in $`T^{}MF^1TN`$, and by $`\overline{R}(X,Y)\xi `$, the curvature tensor of this connection, namely $`(\overline{R}(X,Y)\xi )(Z)`$ $`=R^N(dF(X),dF(Y))\xi (Z)\xi (R^M(X,Y)Z)`$. From Lemma 5.3, for $`p`$ on a neighbourhood of $`p_0`$,
$$d\stackrel{~}{g}_{\mu \overline{\mu }}(\beta )=ig(_\beta dF(\mu ),JdF(\overline{\mu }))ig(_\beta dF(\overline{\mu }),JdF(\mu ))+2\underset{\rho }{}\left(_\beta \mu ,\overline{\rho }\stackrel{~}{g}_{\rho \overline{\mu }}+_\beta \overline{\mu },\rho \stackrel{~}{g}_{\mu \overline{\rho }}\right).$$
Then at $`p_0`$,
$`d\left(d\stackrel{~}{g}_{\mu \overline{\mu }}(\beta )\right)(\overline{\beta })=`$
$`=`$ $`ig(_{\overline{\beta }}(_\beta dF(\mu )),JdF(\overline{\mu }))+ig(_\beta dF(\mu ),_{\overline{\beta }}(JdF(\overline{\mu })\left)\right)`$
$`ig(_{\overline{\beta }}(_\beta dF(\overline{\mu })),JdF(\mu ))ig(_\beta dF(\overline{\mu }),_{\overline{\beta }}(JdF(\mu )\left)\right)`$
$`+2{\displaystyle \underset{\rho }{}}(_{\overline{\beta }}\left(_\beta \mu ,\overline{\rho }\right)\stackrel{~}{g}_{\rho \overline{\mu }}+_{\overline{\beta }}\left(_\beta \overline{\mu },\rho \right)\stackrel{~}{g}_{\mu \overline{\rho }})`$
$`+{\displaystyle \underset{\rho }{}}2_\beta \mu ,\overline{\rho }d\stackrel{~}{g}_{\rho \overline{\mu }}(\overline{\beta })+2_\beta \overline{\mu },\rho d\stackrel{~}{g}_{\mu \overline{\rho }}(\overline{\beta })`$
$`=`$ $`ig(_{\overline{\beta }}(_\beta dF(\mu )),JdF(\overline{\mu }))+ig(_\beta dF(\mu ),J_{\overline{\beta }}dF(\overline{\mu }))`$
$`+ig(_\beta dF(\mu ),JdF(_{\overline{\beta }}\overline{\mu }))ig(_{\overline{\beta }}(_\beta dF(\overline{\mu })),JdF(\mu ))`$
$`ig(_\beta dF(\overline{\mu }),J_{\overline{\beta }}dF(\mu ))ig(_\beta dF(\overline{\mu }),JdF(_{\overline{\beta }}\mu ))`$
$`+\mathrm{cos}\theta _\mu (_{\overline{\beta }}\left(_\beta \mu ,\overline{\mu }\right)+_{\overline{\beta }}(\mu ,_\beta \overline{\mu }\left)\right)+(5.8)`$
$`=`$ $`ig(_{\overline{\beta }}(_\beta dF(\mu )),JdF(\overline{\mu }))`$
$`+ig(_\beta dF(\mu ),J_{\overline{\beta }}dF(\overline{\mu }))+{\displaystyle }_\rho 2i_{\overline{\beta }}\overline{\mu },\rho g_\beta \mu \overline{\rho }+2i_{\overline{\beta }}\overline{\mu },\overline{\rho }g_\beta \mu \rho `$
$`ig(_{\overline{\beta }}(_\beta dF(\overline{\mu })),JdF(\mu ))`$
$`ig(_\beta dF(\overline{\mu }),J_{\overline{\beta }}dF(\mu ))+{\displaystyle }_\rho 2i_{\overline{\beta }}\mu ,\rho g_\beta \overline{\mu }\overline{\rho }2i_{\overline{\beta }}\mu ,\overline{\rho }g_\beta \overline{\mu }\rho `$
$`+\mathrm{cos}\theta _\mu (_{\overline{\beta }}\left(_\beta \mu ,\overline{\mu }\right)+_{\overline{\beta }}(\mu ,_\beta \overline{\mu }\left)\right)`$
$`+(5.8).`$
The term (5.11) vanish because $`_\beta \mu ,\overline{\mu }=\mu ,_\beta \overline{\mu }`$ on a neighbourhood of $`p_0`$. Minimality of $`F`$ implies
$`{\displaystyle \underset{\beta }{}}_{\overline{\beta }}(_\beta dF(\mu ))=`$
$`=`$ $`{\displaystyle \underset{\beta }{}}_{\overline{\beta }}(_\mu dF(\beta ))={\displaystyle }_\beta _{\overline{\beta }}_\mu dF(\beta )+_\mu dF(_{\overline{\beta }}\beta )`$
$`=`$ $`{\displaystyle \underset{\beta }{}}_\mu _{\overline{\beta }}dF(\beta )_{[\mu ,\overline{\beta }]}dF(\beta )+(\overline{R}(\mu ,\overline{\beta })dF)(\beta )+_\mu dF(_{\overline{\beta }}\beta )`$
$`=`$ $`{\displaystyle \underset{\beta }{}}_\mu (_{\overline{\beta }}dF(\beta ))_{\overline{\beta }}dF(_\mu \beta )_{[\mu ,\overline{\beta }]}dF(\beta )`$
$`+R^N(dF(\mu ),dF(\overline{\beta }))dF(\beta )dF(R^M(\mu ,\overline{\beta })\beta )+_\mu dF(_{\overline{\beta }}\beta )`$
$`=`$ $`{\displaystyle \underset{\beta }{}}{\displaystyle \underset{\rho }{}}2_\mu \beta ,\overline{\rho }_{\overline{\beta }}dF(\rho )+{\displaystyle \underset{\rho }{}}2_\mu \beta ,\rho _{\overline{\beta }}dF(\overline{\rho })`$
$`{\displaystyle \underset{\rho }{}}(2_\mu \overline{\beta },\overline{\rho }2_{\overline{\beta }}\mu ,\overline{\rho })_\rho dF(Z_\beta )`$
$`{\displaystyle \underset{\rho }{}}(2_\mu \overline{\beta },\rho 2_{\overline{\beta }}\mu ,\rho )_{\overline{\rho }}dF(Z_\beta )`$
$`+R^N(dF(\mu ),dF(\overline{\beta }))dF(\beta )dF(R^M(\mu ,\overline{\beta })\beta )`$
$`+{\displaystyle \underset{\rho }{}}2_{\overline{\beta }}\beta ,\overline{\rho }_\mu dF(\rho )+{\displaystyle \underset{\rho }{}}2_{\overline{\beta }}\beta ,\rho _\mu dF(\overline{\rho }).`$
Hence
$`(5.9)={\displaystyle \underset{\beta }{}}iR^N(dF(\mu ),dF(\overline{\beta }),dF(\beta ),JdF(\overline{\mu }))\mathrm{cos}\theta _\mu R^M(\mu ,\overline{\beta },\beta ,\overline{\mu })`$
$`+{\displaystyle \underset{\beta \rho }{}}2i_\mu \beta ,\overline{\rho }g_{\overline{\beta }}\rho \overline{\mu }2i_\mu \beta ,\rho g_{\overline{\beta }}\overline{\rho }\overline{\mu }`$
$`+{\displaystyle \underset{\beta \rho }{}}2i(_\mu \overline{\beta },\overline{\rho }+_{\overline{\beta }}\mu ,\overline{\rho })g_\rho \beta \overline{\mu }+2i(_\mu \overline{\beta },\rho +_{\overline{\beta }}\mu ,\rho )g_{\overline{\rho }}\beta \overline{\mu }`$
$`+{\displaystyle \underset{\beta \rho }{}}2i_{\overline{\beta }}\beta ,\overline{\rho }g_\mu \rho \overline{\mu }+2i_{\overline{\beta }}\beta ,\rho g_\mu \overline{\rho }\overline{\mu }.`$
Similarly
$`(5.10)={\displaystyle \underset{\beta }{}}iR^N(dF(\overline{\mu }),dF(\overline{\beta }),dF(\beta ),JdF(\mu ))+\mathrm{cos}\theta _\mu R^M(\overline{\mu },\overline{\beta },\beta ,\mu )`$
$`+{\displaystyle \underset{\beta \rho }{}}2i_{\overline{\mu }}\beta ,\overline{\rho }g_{\overline{\beta }}\rho \mu 2i_{\overline{\mu }}\beta ,\rho g_{\overline{\beta }}\overline{\rho }\mu `$
$`+{\displaystyle \underset{\beta \rho }{}}2i(_{\overline{\mu }}\overline{\beta },\overline{\rho }+_{\overline{\beta }}\overline{\mu },\overline{\rho })g_\rho \beta \mu +2i(_{\overline{\mu }}\overline{\beta },\rho +_{\overline{\beta }}\overline{\mu },\rho )g_{\overline{\rho }}\beta \mu `$
$`+{\displaystyle \underset{\beta \rho }{}}2i_{\overline{\beta }}\beta ,\overline{\rho }g_{\overline{\mu }}\rho \mu +2i_{\overline{\beta }}\beta ,\rho g_{\overline{\mu }}\overline{\rho }\mu .`$
Using Bianchi identity,
$`iR^N(dF(\mu ),dF(\overline{\beta }),dF(\beta ),JdF(\overline{\mu }))iR^N(dF(\overline{\mu }),dF(\overline{\beta }),dF(\beta ),JdF(\mu ))=`$
$`=`$ $`iR^N(dF(\beta ),dF(\mu ),dF(\overline{\beta }),JdF(\overline{\mu }))iR^N(dF(\overline{\beta }),dF(\beta ),dF(\mu ),JdF(\overline{\mu }))`$
$`iR^N(dF(\overline{\mu }),dF(\overline{\beta }),dF(\beta ),JdF(\mu ))`$
$`=`$ $`iR^N(dF(\beta ),dF(\overline{\beta }),dF(\mu ),JdF(\overline{\mu }))+2Im\left(R^N(dF(\beta ),dF(\mu ),dF(\overline{\beta }),JdF(\overline{\mu }))\right),`$
and by Gauss equation, and minimality of $`F`$,
$`{\displaystyle \underset{\beta }{}}R^M(\mu ,\overline{\beta },\beta ,\overline{\mu })R^M(\overline{\mu },\overline{\beta },\beta ,\mu )=`$
$`=`$ $`{\displaystyle \underset{\beta }{}}R^M(\beta ,\mu ,\overline{\beta },\overline{\mu })+R^M(\overline{\beta },\beta ,\mu ,\overline{\mu })R^M(\overline{\mu },\overline{\beta },\beta ,\mu )`$
$`=`$ $`{\displaystyle \underset{\beta }{}}R^M(\beta ,\overline{\beta },\mu ,\overline{\mu })+2R^M(\beta ,\mu ,\overline{\beta },\overline{\mu })`$
$`=`$ $`{\displaystyle \underset{\beta }{}}R^N(dF(\beta ),dF(\overline{\beta }),dF(\mu ),dF(\overline{\mu }))`$
$`g(_\beta dF(\mu ),_{\overline{\beta }}dF(\overline{\mu }))+g(_\beta dF(\overline{\mu }),_{\overline{\beta }}dF(\mu ))`$
$`+2R^N(dF(\beta ),dF(\mu ),dF(\overline{\beta }),dF(\overline{\mu }))`$
$`+2g(_\beta dF(\overline{\beta }),_\mu dF(\overline{\mu }))2g(_\beta dF(\overline{\mu }),_\mu dF(\overline{\beta }))`$
$`=`$ $`{\displaystyle \underset{\beta }{}}R^N(dF(\beta ),dF(\overline{\beta }),dF(\mu ),dF(\overline{\mu }))+2R^N(dF(\beta ),dF(\mu ),dF(\overline{\beta }),dF(\overline{\mu }))`$
$`g(_\beta dF(\mu ),_{\overline{\beta }}dF(\overline{\mu }))g(_\beta dF(\overline{\mu }),_\mu dF(\overline{\beta })).`$
Note that $`R^N(dF(\beta ),dF(\mu ),dF(\overline{\beta }),dF(\overline{\mu }))=Im\left(iR^N(dF(\beta ),dF(\mu ),dF(\overline{\beta }),dF(\overline{\mu }))\right)`$, since it is real. Therefore,
$`{\displaystyle \underset{\beta }{}}d\left(d\stackrel{~}{g}_{\mu \overline{\mu }}(\beta )\right)(\overline{\beta })=`$
$`=`$ $`{\displaystyle \underset{\beta }{}}iR^N(dF(\beta ),dF(\overline{\beta }),dF(\mu ),JdF(\overline{\mu })+i\mathrm{cos}\theta _\mu dF(\overline{\mu }))`$
$`+2Im\left(R^N(dF(\beta ),dF(\mu ),dF(\overline{\beta }),JdF(\overline{\mu })+i\mathrm{cos}\theta _\mu dF(\overline{\mu }))\right)`$
$`\mathrm{cos}\theta _\mu g(_\beta dF(\mu ),_{\overline{\beta }}dF(\overline{\mu }))\mathrm{cos}\theta _\mu g(_\beta dF(\overline{\mu }),_\mu dF(\overline{\beta }))`$
$`+{\displaystyle \underset{\rho }{}}2i_\mu \beta ,\overline{\rho }g_{\overline{\beta }}\rho \overline{\mu }2i_\mu \beta ,\rho g_{\overline{\beta }}\overline{\rho }\overline{\mu }`$
$`+{\displaystyle \underset{\rho }{}}2i(_\mu \overline{\beta },\overline{\rho }+_{\overline{\beta }}\mu ,\overline{\rho })g_\rho \beta \overline{\mu }+2i(_\mu \overline{\beta },\rho +_{\overline{\beta }}\mu ,\rho )g_{\overline{\rho }}\beta \overline{\mu }`$
$`+{\displaystyle \underset{\rho }{}}2i_{\overline{\beta }}\beta ,\overline{\rho }g_\mu \rho \overline{\mu }+2i_{\overline{\beta }}\beta ,\rho g_\mu \overline{\rho }\overline{\mu }`$
$`+{\displaystyle \underset{\rho }{}}2i_{\overline{\mu }}\beta ,\overline{\rho }g_{\overline{\beta }}\rho \mu +2i_{\overline{\mu }}\beta ,\rho g_{\overline{\beta }}\overline{\rho }\mu `$
$`+{\displaystyle \underset{\rho }{}}2i(_{\overline{\mu }}\overline{\beta },\overline{\rho }_{\overline{\beta }}\overline{\mu },\overline{\rho })g_\rho \beta \mu +2(_{\overline{\mu }}\overline{\beta },\rho _{\overline{\beta }}\overline{\mu },\rho )g_{\overline{\rho }}\beta \mu `$
$`+{\displaystyle \underset{\rho }{}}2i_{\overline{\beta }}\beta ,\overline{\rho }g_{\overline{\mu }}\rho \mu 2i_{\overline{\beta }}\beta ,\rho g_{\overline{\mu }}\overline{\rho }\mu `$
$`+ig(_\beta dF(\mu ),J_{\overline{\beta }}dF(\overline{\mu }))ig(_\beta dF(\overline{\mu }),J_{\overline{\beta }}dF(\mu ))`$
$`+{\displaystyle \underset{\rho }{}}2i_{\overline{\beta }}\overline{\mu },\rho g_\beta \mu \overline{\rho }+2i_{\overline{\beta }}\overline{\mu },\overline{\rho }g_\beta \mu \rho `$
$`{\displaystyle \underset{\rho }{}}2i_{\overline{\beta }}\mu ,\rho g_\beta \overline{\mu }\overline{\rho }2i_{\overline{\beta }}\mu ,\overline{\rho }g_\beta \overline{\mu }\rho +(5.8).`$
Using the unitary basis $`\{\frac{\sqrt{2}}{\mathrm{sin}\theta _\rho }\mathrm{\Phi }(\rho ),\frac{\sqrt{2}}{\mathrm{sin}\theta _\rho }\mathrm{\Phi }(\overline{\rho })\}`$ of the normal bundle, and (2.1)
$`(5.12)+(5.15)=`$
$`=`$ $`{\displaystyle \underset{\beta ,\rho }{}}{\displaystyle \frac{2\mathrm{cos}\theta _\mu }{\mathrm{sin}^2\theta _\rho }}\left(|g_\beta \mu \rho |^2+|g_\beta \mu \overline{\rho }|^2\right){\displaystyle \underset{\beta ,\rho }{}}{\displaystyle \frac{2\mathrm{cos}\theta _\mu }{\mathrm{sin}^2\theta _\rho }}\left(|g_\beta \overline{\mu }\rho |^2+|g_\beta \overline{\mu }\overline{\rho }|^2\right)`$
$`{\displaystyle \underset{\beta ,\rho }{}}{\displaystyle \frac{2\mathrm{cos}\theta _\rho }{\mathrm{sin}^2\theta _\rho }}\left(|g_\beta \mu \overline{\rho }|^2|g_\beta \mu \rho |^2\right)+{\displaystyle \underset{\beta ,\rho }{}}{\displaystyle \frac{2\mathrm{cos}\theta _\rho }{\mathrm{sin}^2\theta _\rho }}\left(|g_\beta \overline{\mu }\overline{\rho }|^2|g_\beta \overline{\mu }\rho |^2\right)`$
$`=`$ $`2{\displaystyle \underset{\beta ,\rho }{}}{\displaystyle \frac{(\mathrm{cos}\theta _\rho \mathrm{cos}\theta _\mu )}{\mathrm{sin}^2\theta _\rho }}|g_\beta \mu \rho |^22{\displaystyle \underset{\beta ,\rho }{}}{\displaystyle \frac{(\mathrm{cos}\theta _\rho +\mathrm{cos}\theta _\mu )}{\mathrm{sin}^2\theta _\rho }}|g_\beta \mu \overline{\rho }|^2`$
$`2{\displaystyle \underset{\beta ,\rho }{}}{\displaystyle \frac{(\mathrm{cos}\theta _\rho +\mathrm{cos}\theta _\mu )}{\mathrm{sin}^2\theta _\rho }}|g_\beta \overline{\mu }\rho |^2+2{\displaystyle \underset{\beta ,\rho }{}}{\displaystyle \frac{(\mathrm{cos}\theta _\rho \mathrm{cos}\theta _\mu )}{\mathrm{sin}^2\theta _\rho }}|g_\beta \overline{\mu }\overline{\rho }|^2.`$
Applying lemma 5.3 and since $`_Z\mu ,\overline{\mu }+_Z\overline{\mu },\mu =0`$, we have
$`d\stackrel{~}{g}_{\mu \overline{\mu }}(_{\overline{\beta }}\beta )`$ $`=`$ $`{\displaystyle \underset{\rho }{}}2_{\overline{\beta }}\beta ,\overline{\rho }d\stackrel{~}{g}_{\mu \overline{\mu }}(\rho )+{\displaystyle \underset{\rho }{}}2_{\overline{\beta }}\beta ,\rho d\stackrel{~}{g}_{\mu \overline{\mu }}(\overline{\rho })`$
$`=`$ $`2i{\displaystyle \underset{\rho }{}}\left(_{\overline{\beta }}\beta ,\overline{\rho }g_\rho \mu \overline{\mu }_{\overline{\beta }}\beta ,\overline{\rho }g_\rho \overline{\mu }\mu +_{\overline{\beta }}\beta ,\rho g_{\overline{\rho }}\mu \overline{\mu }_{\overline{\beta }}\beta ,\rho g_{\overline{\rho }}\overline{\mu }\mu \right)`$
$`=`$ $`(5.13)+(5.14).`$
Finally
$`(5.8)`$ $`=`$ $`{\displaystyle \underset{\rho }{}}2_\beta \mu ,\overline{\rho }(ig_{\overline{\beta }}\rho \overline{\mu }ig_{\overline{\beta }}\overline{\mu }\rho )2_\beta \mu ,\overline{\rho }(\mathrm{cos}\theta _\rho \mathrm{cos}\theta _\mu )_{\overline{\beta }}\rho ,\overline{\mu }`$
$`+{\displaystyle \underset{\rho }{}}2_\beta \overline{\mu },\rho (ig_{\overline{\beta }}\mu \overline{\rho }ig_{\overline{\beta }}\overline{\rho }\mu )2_\beta \overline{\mu },\rho (\mathrm{cos}\theta _\mu \mathrm{cos}\theta _\rho )_{\overline{\beta }}\mu ,\overline{\rho }`$
$`=`$ $`{\displaystyle \underset{\rho }{}}2i_\beta \mu ,\overline{\rho }g_{\overline{\beta }}\rho \overline{\mu }2i_\beta \mu ,\overline{\rho }g_{\overline{\beta }}\overline{\mu }\rho +2i_\beta \overline{\mu },\rho g_{\overline{\beta }}\mu \overline{\rho }2i_\beta \overline{\mu },\rho g_{\overline{\beta }}\overline{\rho }\mu `$
$`2{\displaystyle \underset{\rho }{}}(\mathrm{cos}\theta _\mu \mathrm{cos}\theta _\rho )\left(|_\beta \mu ,\overline{\rho }|^2+|_{\overline{\beta }}\mu ,\overline{\rho }|^2\right).`$
These expressions lead to the expression of the lemma. $`\mathbf{}`$
Finally, we have
###### Proposition 5.1
If F is minimal without complex directions, then for each $`0k2n`$ at each $`p_0\mathrm{\Omega }_{2k}^0`$,
$`\mathrm{}\kappa `$ $`=`$ $`4i{\displaystyle \underset{\beta }{}}Ricci^N(JdF(\beta ),dF(\overline{\beta }))`$
$`+{\displaystyle \underset{\beta ,\mu }{}}{\displaystyle \frac{32}{\mathrm{sin}^2\theta _\mu }}Im\left(R^N(dF(\beta ),dF(\mu ),dF(\overline{\beta }),JdF(\overline{\mu })+i\mathrm{cos}\theta _\mu dF(\overline{\mu }))\right)`$
$`{\displaystyle \underset{\beta ,\mu ,\rho }{}}{\displaystyle \frac{64(\mathrm{cos}\theta _\mu +\mathrm{cos}\theta _\rho )}{\mathrm{sin}^2\theta _\mu \mathrm{sin}^2\theta _\rho }}Re\left(g_\beta \mu \overline{\rho }g_{\overline{\beta }}\rho \overline{\mu }\right)`$
$`+{\displaystyle \underset{\beta ,\mu ,\rho }{}}{\displaystyle \frac{32(\mathrm{cos}\theta _\rho \mathrm{cos}\theta _\mu )}{\mathrm{sin}^2\theta _\mu \mathrm{sin}^2\theta _\rho }}(|g_\beta \mu \rho |^2+|g_{\overline{\beta }}\mu \rho |^2)`$
$`+{\displaystyle \underset{\beta ,\mu ,\rho }{}}{\displaystyle \frac{32(\mathrm{cos}\theta _\mu +\mathrm{cos}\theta _\rho )}{\mathrm{sin}^2\theta _\mu }}\left(|_\beta \mu ,\rho |^2+|_{\overline{\beta }}\mu ,\rho |^2\right).`$
Proof. From (5.7) and Lemma 5.4 we get
$`2\mathrm{}\kappa =`$
$`=`$ $`+{\displaystyle \underset{\beta ,\mu ,\rho }{}}{\displaystyle \frac{64(\mathrm{cos}\theta _\mu +\mathrm{cos}\theta _\rho )}{\mathrm{sin}^2\theta _\mu \mathrm{sin}^2\theta _\rho }}d\stackrel{~}{g}_{\mu \overline{\rho }}(\overline{\beta })d\stackrel{~}{g}_{\rho \overline{\mu }}(\beta )`$
$`+{\displaystyle \underset{\beta ,\mu }{}}{\displaystyle \frac{32i}{\mathrm{sin}^2\theta _\mu }}R^N(dF(\beta ),dF(\overline{\beta }),dF(\mu ),JdF(\overline{\mu })+i\mathrm{cos}\theta _\mu dF(\overline{\mu }))`$
$`+{\displaystyle \underset{\beta ,\mu }{}}{\displaystyle \frac{64}{\mathrm{sin}^2\theta _\mu }}Im\left(R^N(dF(\beta ),dF(\mu ),dF(\overline{\beta }),JdF(\overline{\mu })+i\mathrm{cos}\theta _\mu dF(\overline{\mu }))\right)`$
$`+{\displaystyle \underset{\beta ,\mu ,\rho }{}}{\displaystyle \frac{64(\mathrm{cos}\theta _\rho \mathrm{cos}\theta _\mu )}{\mathrm{sin}^2\theta _\mu \mathrm{sin}^2\theta _\rho }}(|g_\beta \mu \rho |^2+|g_\beta \overline{\mu }\overline{\rho }|^2)`$
$`{\displaystyle \underset{\beta ,\mu ,\rho }{}}{\displaystyle \frac{64(\mathrm{cos}\theta _\rho +\mathrm{cos}\theta _\mu )}{\mathrm{sin}^2\theta _\mu \mathrm{sin}^2\theta _\rho }}(|g_\beta \mu \overline{\rho }|^2+|g_\beta \overline{\mu }\rho |^2)`$
$`+{\displaystyle \underset{\beta ,\mu ,\rho }{}}{\displaystyle \frac{64i}{\mathrm{sin}^2\theta _\mu }}_\mu \beta ,\overline{\rho }g_{\overline{\beta }}\rho \overline{\mu }{\displaystyle \frac{64i}{\mathrm{sin}^2\theta _\mu }}_\mu \beta ,\rho g_{\overline{\beta }}\overline{\rho }\overline{\mu }{\displaystyle \frac{64i}{\mathrm{sin}^2\theta _\mu }}_\mu \overline{\beta },\overline{\rho }g_\rho \beta \overline{\mu }`$
$`+{\displaystyle \underset{\beta ,\mu ,\rho }{}}{\displaystyle \frac{64i}{\mathrm{sin}^2\theta _\mu }}_{\overline{\beta }}\mu ,\overline{\rho }g_\beta \rho \overline{\mu }{\displaystyle \frac{64i}{\mathrm{sin}^2\theta _\mu }}_\mu \overline{\beta },\rho g_{\overline{\rho }}\beta \overline{\mu }+{\displaystyle \frac{64i}{\mathrm{sin}^2\theta _\mu }}_{\overline{\beta }}\mu ,\rho g_{\overline{\rho }}\beta \overline{\mu }`$
$`+{\displaystyle \underset{\beta ,\mu ,\rho }{}}{\displaystyle \frac{64i}{\mathrm{sin}^2\theta _\mu }}_{\overline{\mu }}\beta ,\overline{\rho }g_{\overline{\beta }}\rho \mu +{\displaystyle \frac{64i}{\mathrm{sin}^2\theta _\mu }}_{\overline{\mu }}\beta ,\rho g_{\overline{\beta }}\overline{\rho }\mu +{\displaystyle \frac{64i}{\mathrm{sin}^2\theta _\mu }}_{\overline{\mu }}\overline{\beta },\overline{\rho }g_\rho \beta \mu `$
$`+{\displaystyle \underset{\beta ,\mu ,\rho }{}}{\displaystyle \frac{64i}{\mathrm{sin}^2\theta _\mu }}_{\overline{\beta }}\overline{\mu },\overline{\rho }g_\rho \beta \mu +{\displaystyle \frac{64i}{\mathrm{sin}^2\theta _\mu }}_{\overline{\mu }}\overline{\beta },\rho g_{\overline{\rho }}\beta \mu {\displaystyle \frac{64i}{\mathrm{sin}^2\theta _\mu }}_{\overline{\beta }}\overline{\mu },\rho g_{\overline{\rho }}\beta \mu `$
$`+{\displaystyle \underset{\beta ,\mu ,\rho }{}}{\displaystyle \frac{64i}{\mathrm{sin}^2\theta _\mu }}_{\overline{\beta }}\overline{\mu },\overline{\rho }g_\beta \mu \rho +{\displaystyle \frac{64i}{\mathrm{sin}^2\theta _\mu }}_{\overline{\beta }}\overline{\mu },\rho g_\beta \mu \overline{\rho }{\displaystyle \frac{64i}{\mathrm{sin}^2\theta _\mu }}_{\overline{\beta }}\mu ,\rho g_\beta \overline{\mu }\overline{\rho }`$
$`+{\displaystyle \underset{\beta ,\mu ,\rho }{}}{\displaystyle \frac{64i}{\mathrm{sin}^2\theta _\mu }}_{\overline{\beta }}\mu ,\overline{\rho }g_\beta \overline{\mu }\rho +{\displaystyle \frac{64i}{\mathrm{sin}^2\theta _\mu }}_\beta \mu ,\overline{\rho }g_{\overline{\beta }}\rho \overline{\mu }{\displaystyle \frac{64i}{\mathrm{sin}^2\theta _\mu }}_\beta \mu ,\overline{\rho }g_{\overline{\beta }}\overline{\mu }\rho `$
$`+{\displaystyle \underset{\beta ,\mu ,\rho }{}}{\displaystyle \frac{64i}{\mathrm{sin}^2\theta _\mu }}_\beta \overline{\mu },\rho g_{\overline{\beta }}\mu \overline{\rho }{\displaystyle \frac{64i}{\mathrm{sin}^2\theta _\mu }}_\beta \overline{\mu },\rho g_{\overline{\beta }}\overline{\rho }\mu `$
$`{\displaystyle \underset{\beta ,\mu ,\rho }{}}{\displaystyle \frac{64(\mathrm{cos}\theta _\mu \mathrm{cos}\theta _\rho )}{\mathrm{sin}^2\theta _\mu }}\left(|_\beta \mu ,\overline{\rho }|^2+|_{\overline{\beta }}\mu ,\overline{\rho }|^2\right).`$
Interchanging $`\rho `$ with $`\beta `$ in the first term of (5.16) (that we named by (5.16)(1), and similarly to other equations), we see that $`(5.16)(1)+(5.17)(2)=0`$. Interchanging $`\rho `$ with $`\beta `$ in (5.18)(1), we get $`(5.18)(1)+(5.19)(2)=0`$. In (5.16)(2), $`_\mu \beta ,\rho `$ is skew-symmetric on $`\rho `$ and $`\beta `$, and $`g_{\overline{\beta }}\overline{\rho }\overline{\mu }`$ is symmetric on $`\rho `$ and $`\beta `$. Hence $`(5.16)(2)=0`$. Similarly $`(5.16)(3)=(5.18)(2)=(5.18)(3)=0`$. If we interchange $`\rho `$ with $`\mu `$ in (5.17)(1),
$$(5.17)(1)+(5.20)(2)=\underset{\beta ,\mu ,\rho }{}\frac{64i(\mathrm{sin}^2\theta _\mu \mathrm{sin}^2\theta _\rho )}{\mathrm{sin}^2\theta _\mu \mathrm{sin}^2\theta _\rho }_{\overline{\beta }}\overline{\mu },\rho g_\beta \mu \overline{\rho }.$$
Interchanging $`\rho `$ with $`\mu `$ in (5.17)(3), we get
$$(5.17)(3)+(5.20)(3)=\underset{\beta ,\mu ,\rho }{}\frac{64i(\mathrm{sin}^2\theta _\mu +\mathrm{sin}^2\theta _\rho )}{\mathrm{sin}^2\theta _\mu \mathrm{sin}^2\theta _\rho }_{\overline{\beta }}\mu ,\rho g_\beta \overline{\mu }\overline{\rho }.$$
Interchanging $`\rho `$ with $`\mu `$ in (5.19)(1), we get
$$(5.19)(1)+(5.20)(1)=\underset{\beta ,\mu ,\rho }{}\frac{64i(\mathrm{sin}^2\theta _\mu +\mathrm{sin}^2\theta _\rho )}{\mathrm{sin}^2\theta _\mu \mathrm{sin}^2\theta _\rho }_{\overline{\beta }}\overline{\mu },\overline{\rho }g_\beta \mu \rho .$$
Interchanging $`\rho `$ with $`\mu `$ in (5.19)(3), we get
$$(5.19)(3)+(5.21)(1)=\underset{\beta ,\mu ,\rho }{}\frac{64i(\mathrm{sin}^2\theta _\mu \mathrm{sin}^2\theta _\rho )}{\mathrm{sin}^2\theta _\mu \mathrm{sin}^2\theta _\rho }_{\overline{\beta }}\mu ,\overline{\rho }g_\beta \overline{\mu }\rho .$$
Interchanging $`\rho `$ with $`\mu `$ in (5.21)(2),
$$(5.21)(2)+(5.22)(1)=\underset{\beta ,\mu ,\rho }{}\frac{64i(\mathrm{sin}^2\theta _\mu +\mathrm{sin}^2\theta _\rho )}{\mathrm{sin}^2\theta _\mu \mathrm{sin}^2\theta _\rho }_\beta \overline{\mu },\rho g_{\overline{\beta }}\mu \overline{\rho }.$$
Interchanging $`\rho `$ with $`\mu `$ in (5.22)(2), we obtain
$$(5.22)(2)+(5.21)(3)=\underset{\beta ,\mu ,\rho }{}\frac{64i(\mathrm{sin}^2\theta _\mu \mathrm{sin}^2\theta _\rho )}{\mathrm{sin}^2\theta _\mu \mathrm{sin}^2\theta _\rho }_\beta \mu ,\overline{\rho }g_{\overline{\beta }}\overline{\mu }\rho .$$
Therefore,
$`2\mathrm{}\kappa =`$ (5.37)
$`=`$ $`{\displaystyle \underset{\beta ,\mu ,\rho }{}}{\displaystyle \frac{64(\mathrm{cos}\theta _\mu +\mathrm{cos}\theta _\rho )}{\mathrm{sin}^2\theta _\mu \mathrm{sin}^2\theta _\rho }}d\stackrel{~}{g}_{\mu \overline{\rho }}(\overline{\beta })d\stackrel{~}{g}_{\rho \overline{\mu }}(\beta )`$
$`+{\displaystyle \underset{\beta ,\mu }{}}{\displaystyle \frac{32i}{\mathrm{sin}^2\theta _\mu }}R^N(dF(\beta ),dF(\overline{\beta }),dF(\mu ),JdF(\overline{\mu })+i\mathrm{cos}\theta _\mu dF(\overline{\mu }))`$
$`+{\displaystyle \underset{\beta ,\mu }{}}{\displaystyle \frac{64}{\mathrm{sin}^2\theta _\mu }}Im\left(R^N(dF(\beta ),dF(\mu ),dF(\overline{\beta }),JdF(\overline{\mu })+i\mathrm{cos}\theta _\mu dF(\overline{\mu }))\right)`$
$`+{\displaystyle \underset{\beta ,\mu ,\rho }{}}{\displaystyle \frac{64(\mathrm{cos}\theta _\rho \mathrm{cos}\theta _\mu )}{\mathrm{sin}^2\theta _\mu \mathrm{sin}^2\theta _\rho }}|g_\beta \mu \rho |^2`$
$`{\displaystyle \underset{\beta ,\mu ,\rho }{}}{\displaystyle \frac{64(\mathrm{cos}\theta _\rho +\mathrm{cos}\theta _\mu )}{\mathrm{sin}^2\theta _\mu \mathrm{sin}^2\theta _\rho }}|g_\beta \mu \overline{\rho }|^2`$
$`{\displaystyle \underset{\beta ,\mu ,\rho }{}}{\displaystyle \frac{64(\mathrm{cos}\theta _\rho +\mathrm{cos}\theta _\mu )}{\mathrm{sin}^2\theta _\mu \mathrm{sin}^2\theta _\rho }}|g_\beta \overline{\mu }\rho |^2`$
$`+{\displaystyle \underset{\beta ,\mu ,\rho }{}}{\displaystyle \frac{64(\mathrm{cos}\theta _\rho \mathrm{cos}\theta _\mu )}{\mathrm{sin}^2\theta _\mu \mathrm{sin}^2\theta _\rho }}|g_\beta \overline{\mu }\overline{\rho }|^2`$
$`{\displaystyle \underset{\beta ,\mu ,\rho }{}}{\displaystyle \frac{64i(\mathrm{sin}^2\theta _\mu \mathrm{sin}^2\theta _\rho )}{\mathrm{sin}^2\theta _\mu \mathrm{sin}^2\theta _\rho }}_{\overline{\beta }}\overline{\mu },\rho g_\beta \mu \overline{\rho }`$
$`{\displaystyle \underset{\beta ,\mu ,\rho }{}}{\displaystyle \frac{64i(\mathrm{sin}^2\theta _\mu +\mathrm{sin}^2\theta _\rho )}{\mathrm{sin}^2\theta _\mu \mathrm{sin}^2\theta _\rho }}_{\overline{\beta }}\mu ,\rho g_\beta \overline{\mu }\overline{\rho }`$
$`+{\displaystyle \underset{\beta ,\mu ,\rho }{}}{\displaystyle \frac{64i(\mathrm{sin}^2\theta _\mu +\mathrm{sin}^2\theta _\rho )}{\mathrm{sin}^2\theta _\mu \mathrm{sin}^2\theta _\rho }}_{\overline{\beta }}\overline{\mu },\overline{\rho }g_\beta \mu \rho `$
$`+{\displaystyle \underset{\beta ,\mu ,\rho }{}}{\displaystyle \frac{64i(\mathrm{sin}^2\theta _\mu \mathrm{sin}^2\theta _\rho )}{\mathrm{sin}^2\theta _\mu \mathrm{sin}^2\theta _\rho }}_{\overline{\beta }}\mu ,\overline{\rho }g_\beta \overline{\mu }\rho `$
$`+{\displaystyle \underset{\beta ,\mu ,\rho }{}}{\displaystyle \frac{64i(\mathrm{sin}^2\theta _\mu +\mathrm{sin}^2\theta _\rho )}{\mathrm{sin}^2\theta _\mu \mathrm{sin}^2\theta _\rho }}_\beta \overline{\mu },\rho g_{\overline{\beta }}\mu \overline{\rho }`$
$`+{\displaystyle \underset{\beta ,\mu ,\rho }{}}{\displaystyle \frac{64i(\mathrm{sin}^2\theta _\mu \mathrm{sin}^2\theta _\rho )}{\mathrm{sin}^2\theta _\mu \mathrm{sin}^2\theta _\rho }}_\beta \mu ,\overline{\rho }g_{\overline{\beta }}\overline{\mu }\rho `$
$`{\displaystyle \underset{\beta ,\mu ,\rho }{}}{\displaystyle \frac{64(\mathrm{cos}\theta _\mu \mathrm{cos}\theta _\rho )}{\mathrm{sin}^2\theta _\mu }}\left(|_\beta \mu ,\overline{\rho }|^2+|_{\overline{\beta }}\mu ,\overline{\rho }|^2\right).`$
By Lemma 5.3,
$`(5.23)`$ $`=`$ $`{\displaystyle \underset{\beta ,\mu ,\rho }{}}{\displaystyle \frac{64(\mathrm{cos}\theta _\mu +\mathrm{cos}\theta _\rho )}{\mathrm{sin}^2\theta _\mu \mathrm{sin}^2\theta _\rho }}(ig_{\overline{\beta }}\mu \overline{\rho }ig_{\overline{\beta }}\overline{\rho }\mu (\mathrm{cos}\theta _\mu \mathrm{cos}\theta _\rho )_{\overline{\beta }}\mu ,\overline{\rho })`$ (5.44)
$`\left(ig_\beta \rho \overline{\mu }ig_\beta \overline{\mu }\rho (\mathrm{cos}\theta _\rho \mathrm{cos}\theta _\mu )_\beta \rho ,\overline{\mu }\right)`$
$`=`$ $`{\displaystyle \underset{\beta ,\mu ,\rho }{}}{\displaystyle \frac{64(\mathrm{cos}\theta _\mu +\mathrm{cos}\theta _\rho )}{\mathrm{sin}^2\theta _\mu \mathrm{sin}^2\theta _\rho }}g_{\overline{\beta }}\mu \overline{\rho }g_\beta \rho \overline{\mu }`$
$`+{\displaystyle \underset{\beta ,\mu ,\rho }{}}{\displaystyle \frac{64(\mathrm{cos}\theta _\mu +\mathrm{cos}\theta _\rho )}{\mathrm{sin}^2\theta _\mu \mathrm{sin}^2\theta _\rho }}|g_\beta \overline{\mu }\rho |^2`$
$`+{\displaystyle \underset{\beta ,\mu ,\rho }{}}{\displaystyle \frac{64i(\mathrm{cos}^2\theta _\mu \mathrm{cos}^2\theta _\rho )}{\mathrm{sin}^2\theta _\mu \mathrm{sin}^2\theta _\rho }}g_{\overline{\beta }}\mu \overline{\rho }_\beta \rho ,\overline{\mu }`$
$`+{\displaystyle \underset{\beta ,\mu ,\rho }{}}{\displaystyle \frac{64(\mathrm{cos}\theta _\mu +\mathrm{cos}\theta _\rho )}{\mathrm{sin}^2\theta _\mu \mathrm{sin}^2\theta _\rho }}|g_\beta \rho \overline{\mu }|^2`$
$`{\displaystyle \underset{\beta ,\mu ,\rho }{}}{\displaystyle \frac{64(\mathrm{cos}\theta _\mu +\mathrm{cos}\theta _\rho )}{\mathrm{sin}^2\theta _\mu \mathrm{sin}^2\theta _\rho }}g_\beta \overline{\mu }\rho g_{\overline{\beta }}\overline{\rho }\mu `$
$`{\displaystyle \underset{\beta ,\mu ,\rho }{}}{\displaystyle \frac{64i(\mathrm{cos}^2\theta _\mu \mathrm{cos}^2\theta _\rho )}{\mathrm{sin}^2\theta _\mu \mathrm{sin}^2\theta _\rho }}_\beta \rho ,\overline{\mu }g_{\overline{\beta }}\overline{\rho }\mu `$
$`{\displaystyle \underset{\beta ,\mu ,\rho }{}}{\displaystyle \frac{64i(\mathrm{cos}^2\theta _\mu \mathrm{cos}^2\theta _\rho )}{\mathrm{sin}^2\theta _\mu \mathrm{sin}^2\theta _\rho }}_{\overline{\beta }}\mu ,\overline{\rho }g_\beta \rho \overline{\mu }`$
$`+{\displaystyle \underset{\beta ,\mu ,\rho }{}}{\displaystyle \frac{64i(\mathrm{cos}^2\theta _\mu \mathrm{cos}^2\theta _\rho )}{\mathrm{sin}^2\theta _\mu \mathrm{sin}^2\theta _\rho }}_{\overline{\beta }}\mu ,\overline{\rho }g_\beta \overline{\mu }\rho `$
$`+{\displaystyle \underset{\beta ,\mu ,\rho }{}}{\displaystyle \frac{64(\mathrm{cos}^2\theta _\mu \mathrm{cos}^2\theta _\rho )}{\mathrm{sin}^2\theta _\mu \mathrm{sin}^2\theta _\rho }}(\mathrm{cos}\theta _\rho \mathrm{cos}\theta _\mu )_{\overline{\beta }}\mu ,\overline{\rho }_\beta \rho ,\overline{\mu }.`$
Immediately we have, $`(5.27)+(5.36)=(5.32)+(5.41)=(5.33)+(5.37)=0`$, and interchanging $`\mu `$ with $`\rho `$ in (5.26), (5.34) and in (5.40), we get, $`(5.26)+(5.38)=(5.29)+(5.40)=(5.34)+(5.39)=0`$. Note that
$$\underset{\mu ,\rho }{}\frac{(\mathrm{cos}\theta _\mu \mathrm{cos}\theta _\rho )}{\mathrm{sin}^2\theta _\mu }|_\beta \mu ,\overline{\rho }|^2=\underset{\mu ,\rho }{}\frac{(\mathrm{cos}\theta _\rho \mathrm{cos}\theta _\mu )}{\mathrm{sin}^2\theta _\rho }|_{\overline{\beta }}\mu ,\overline{\rho }|^2.$$
Hence $`(5.35)+(5.42)=0`$. Then,
$`2\mathrm{}\kappa `$ $`=`$ $`{\displaystyle \underset{\beta ,\mu }{}}{\displaystyle \frac{32i}{\mathrm{sin}^2\theta _\mu }}R^N(dF(\beta ),dF(\overline{\beta }),dF(\mu ),JdF(\overline{\mu })+i\mathrm{cos}\theta _\mu dF(\overline{\mu }))`$ (5.48)
$`+{\displaystyle \underset{\beta ,\mu }{}}{\displaystyle \frac{64}{\mathrm{sin}^2\theta _\mu }}Im\left(R^N(dF(\beta ),dF(\mu ),dF(\overline{\beta }),JdF(\overline{\mu })+i\mathrm{cos}\theta _\mu dF(\overline{\mu }))\right)`$
$`+{\displaystyle \underset{\beta ,\mu ,\rho }{}}{\displaystyle \frac{64(\mathrm{cos}\theta _\mu +\mathrm{cos}\theta _\rho )}{\mathrm{sin}^2\theta _\mu \mathrm{sin}^2\theta _\rho }}(g_{\overline{\beta }}\mu \overline{\rho }g_\beta \rho \overline{\mu }+g_\beta \overline{\mu }\rho g_{\overline{\beta }}\overline{\rho }\mu )`$
$`+{\displaystyle \underset{\beta ,\mu ,\rho }{}}{\displaystyle \frac{64(\mathrm{cos}\theta _\rho \mathrm{cos}\theta _\mu )}{\mathrm{sin}^2\theta _\mu \mathrm{sin}^2\theta _\rho }}(|g_\beta \mu \rho |^2+|g_{\overline{\beta }}\mu \rho |^2)`$
$`{\displaystyle \underset{\beta ,\mu ,\rho }{}}{\displaystyle \frac{64i(\mathrm{sin}^2\theta _\mu +\mathrm{sin}^2\theta _\rho )}{\mathrm{sin}^2\theta _\mu \mathrm{sin}^2\theta _\rho }}_{\overline{\beta }}\mu ,\rho g_\beta \overline{\mu }\overline{\rho }`$
$`+{\displaystyle \underset{\beta ,\mu ,\rho }{}}{\displaystyle \frac{64i(\mathrm{sin}^2\theta _\mu +\mathrm{sin}^2\theta _\rho )}{\mathrm{sin}^2\theta _\mu \mathrm{sin}^2\theta _\rho }}_{\overline{\beta }}\overline{\mu },\overline{\rho }g_\beta \mu \rho .`$
Using Lemma 5.3, and interchanging $`\rho `$ by $`\mu `$ when necessary,
$`(5.45)+(5.46)=`$
$`=`$ $`{\displaystyle \underset{\beta ,\mu ,\rho }{}}{\displaystyle \frac{64i}{\mathrm{sin}^2\theta _\rho }}_{\overline{\beta }}\mu ,\rho g_\beta \overline{\mu }\overline{\rho }{\displaystyle \frac{64i}{\mathrm{sin}^2\theta _\mu }}_{\overline{\beta }}\mu ,\rho g_\beta \overline{\mu }\overline{\rho }+{\displaystyle \frac{64i}{\mathrm{sin}^2\theta _\mu }}_{\overline{\beta }}\overline{\mu },\overline{\rho }g_\beta \mu \rho +{\displaystyle \frac{64i}{\mathrm{sin}^2\theta _\rho }}_{\overline{\beta }}\overline{\mu },\overline{\rho }g_\beta \mu \rho `$
$`=`$ $`{\displaystyle \underset{\beta ,\mu ,\rho }{}}{\displaystyle \frac{64i}{\mathrm{sin}^2\theta _\mu }}_{\overline{\beta }}\mu ,\rho \left(g_\beta \overline{\mu }\overline{\rho }g_\beta \overline{\rho }\overline{\mu }\right)+{\displaystyle \underset{\beta ,\mu ,\rho }{}}{\displaystyle \frac{64i}{\mathrm{sin}^2\theta _\mu }}_{\overline{\beta }}\overline{\mu },\overline{\rho }\left(g_\beta \mu \rho g_\beta \rho \mu \right)`$
$`=`$ $`{\displaystyle \underset{\beta ,\mu ,\rho }{}}{\displaystyle \frac{64}{\mathrm{sin}^2\theta _\mu }}_{\overline{\beta }}\mu ,\rho (\mathrm{cos}\theta _\mu +\mathrm{cos}\theta _\rho )_\beta \overline{\mu },\overline{\rho }+{\displaystyle \frac{64}{\mathrm{sin}^2\theta _\mu }}_{\overline{\beta }}\overline{\mu },\overline{\rho }(\mathrm{cos}\theta _\mu +\mathrm{cos}\theta _\rho )_\beta \mu ,\rho `$
$`=`$ $`{\displaystyle \underset{\beta ,\mu ,\rho }{}}{\displaystyle \frac{64(\mathrm{cos}\theta _\mu +\mathrm{cos}\theta _\rho )}{\mathrm{sin}^2\theta _\mu }}\left(|_\beta \mu ,\rho |^2+|_{\overline{\beta }}\mu ,\rho |^2\right).`$
Obviously
$$(5.44)=\underset{\beta ,\mu ,\rho }{}\frac{128(\mathrm{cos}\theta _\mu +\mathrm{cos}\theta _\rho )}{\mathrm{sin}^2\theta _\mu \mathrm{sin}^2\theta _\rho }Re\left(g_\beta \mu \overline{\rho }g_{\overline{\beta }}\rho \overline{\mu }\right).$$
From (1.4), (2.1), and the $`J`$-invariance of $`Ricci`$, $`(5.43)=8i_\beta Ricci^N(JdF(\beta ),dF(\overline{\beta }))`$, and the expression of the Proposition follows.
Acknowledgments
We would like to thank very much Professor James Eells for helpful discussions and encouragement.
References
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\[Ch-W\] S.S. Chern & J.G. Wolfson, Minimal surfaces by moving frames, Amer. J. Math. 105 (1983), 59-83.
\[E-L\] J. Eells & L. Lemaire, Selected topics in harmonic maps, C.B.M.S. Regional Conf. Series 50, A.M.S. (1983).
\[Mi-Mo\] M.J. Micallef & J.D. Moore, Minimal two-spheres and the topology of manifolds with positive curvature on totally isotropic two-planes, Annals of Math. 127 (1988), 199-227
\[O\] Y. Ohnita, Minimal surfaces with constant curvature and Kähler angle in complex space forms, Tsukuba J. Math. 13 No1 (1989), 191-207.
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\[S-V\] I. Salavessa & G. Valli, Broadly-Pluriminimal Submanifolds of Kähler-Einstein Manifolds, preprint submitted for publication.
\[W\] J. G. Wolfson, Minimal Surfaces in Kähler Surfaces and Ricci Curvature, J. Diff. Geom, 29 (1989), 281–294. |
warning/0002/hep-ph0002109.html | ar5iv | text | # Neutrinoless double beta decay and new physics in the neutrino sector
## 1 Introduction
Double beta decay corresponds to two single beta decays occuring in one nucleus and converts a nucleus (Z,A) into a nucleus (Z+2,A). While even the standard model (SM) allowed process emitting two antineutrinos
$${}_{Z}{}^{A}X_{Z+2}^AX+2e^{}+2\overline{\nu }_e$$
(1)
is one of the rarest processes in nature with half lives in the region of $`10^{2124}`$ years, more interesting is the search for the neutrinoless mode ($`0\nu \beta \beta `$),
$${}_{Z}{}^{A}X_{Z+2}^AX+2e^{}$$
(2)
which violates lepton number by two units and thus implies physics beyond the SM .
The most sensitive experiment so far, the Heidelberg–Moscow experiment is searching for the $`0\nu \beta \beta `$ decay of <sup>76</sup>Ge. The results after 31 kg y measuring time using digital pulse shape analysis correspond to a conservative half life limit of
$`T_{1/2}^{0\nu \beta \beta }`$ $`>`$ $`1.810^{25}y(90\%C.L.),`$
$`T_{1/2}^{0\nu \beta \beta }`$ $`>`$ $`3.010^{25}y(68\%C.L.).`$ (3)
To render possible a further breakthrough in search for neutrino masses and physics beyond the SM, GENIUS, an experiment operating a large amount of naked Ge–detectors in a liquid nitrogen shielding, has been proposed . Operating 288 enriched <sup>76</sup>Ge detectors with a total mass of 1 ton inside a nitrogen tank of $``$ 12 m height and diameter, one could access half lifes of $`T_{1/2}^{0\nu \beta \beta }=610^{27}y`$ after one year of measurement. A ten ton version would reach a final sensitivity of $`T_{1/2}^{0\nu \beta \beta }=610^{29}y`$ within 10 years of measurement time.
## 2 Neutrino masses and oscillations
The search for $`0\nu \beta \beta `$ decay exchanging a massive left–handed Majorana neutrino between two SM vertices (contribution a) in fig. 1) at present provides the most sensitive approach to determine an absolute neutrino mass and also a unique possibility to distinguish between the Dirac or Majorana nature of the neutrino. With the recent half life limit of the Heidelberg–Moscow experiment the following conservative limit on the effective neutrino mass
$`m_\nu ={\displaystyle \underset{i}{}}U_{ei}^2m_i0.36eV(90\%C.L.)`$
$`m_\nu ={\displaystyle \underset{i}{}}U_{ei}^2m_i0.28eV(68\%C.L.)`$ (4)
can be deduced. Here the sum extends over light mass eigenstates $`m_i`$ only. The GENIUS project could access effective neutrino masses down to $`10^2`$ eV or even $`10^3`$ eV in the 1 ton or 10 ton version, respectively. Since in specific models of neutrino masses the quantity $`m`$ can be related to the oscillation parameters $`\mathrm{\Delta }m^2`$ and $`\mathrm{sin}^22\theta `$, these bounds imply restrictive bounds on the neutrino mass spectrum, which in most cases are more stringent than the bounds from precision measurements of the CMB by MAP and Planck. The extreme cases discussed here are degenerate and hierarchical models (for a detailed discussion see ).
Degenerate models: Such models have been proposed to get a large mass scale for neutrinos acting as hot dark matter and at the same time accomodate two of the three neutrino anomalies (solar, atmospheric & LSND neutrinos). In this framework neutrino masses up to (few) eV are predicted. Thus large mixing (vacuum or MSW LMA oscillations) and strong cancellations are required to be consistent with eq. 4. E.g., maximal cancellation, the MSW LMA bestfit and a $`\mathrm{\Lambda }`$CHDM model with a Hubble constant of $`h0.5`$ implies a value of $`m=0.15`$ eV just below the present limit.
Hierarchical models: Such models give less optimistic predictions for $`0\nu \beta \beta `$ decay. However, assuming the MSW LMA solution, still sizable contributions up to $`m=`$ (few) $`10^2`$ in the reach of GENIUS are possible. Since $`0\nu \beta \beta `$ decay is most sensitive in the large $`\mathrm{\Delta }m^2`$ region of the LMA solution, it may provide complementary informations to the search for day-night effects in solar neutrinos.
For a super-heavy left-handed neutrino a bound of $`m_H=\left(_j\frac{U_{ej}^2}{m_j}\right)^1>910^7GeV`$ can be deduced, with heavy mass eigenstates $`m_j`$. This constraint makes a heavy neutrino unobservable at linear colliders except in the most contrived scenarios .
## 3 New interactions
Besides the exchange of massive Majorana neutrinos between two SM vertices, a variety of theories beyond the SM predicts new lepton number violating interactions contributing to neutrinoless double beta decay, leading to the idea to construct the general double beta decay rate allowed by Lorentz–invariance . This approach allows to constrain lepton number violating parameters in arbitrary models.
For the long range part of the decay rate with two separable vertices and light neutrino exchange in between (contributions b) and c) in fig. 1), one has to consider the Lorentz-invariant contractions of six projections with defined helicity both for the leptonic ($`j_\alpha `$) and hadronic ($`J_\alpha `$) current. The general Lagrangian can be written in terms of effective couplings $`ϵ_\beta ^\alpha `$, which correspond to the pointlike vertices at the Fermi scale so that Fierz rearrangement is applicable:
$$=\frac{G_F}{\sqrt{2}}\{j_{VA}^\mu J_{VA,\mu }^{}+\underset{\alpha ,\beta }{\overset{^{}}{}}ϵ_\alpha ^\beta j_\beta J_\alpha ^{}\}$$
(5)
with the combinations of hadronic and leptonic Lorentz currents of defined helicity $`\alpha ,\beta =VA,V+A,SP,S+P,T_L,T_R`$. The prime indicates that the sum runs over all contractions allowed by Lorentz–invariance, except for $`\alpha =\beta =VA`$. Here $`ϵ_\alpha ^\beta `$ denotes the strength of the non–SM couplings. For the helicity suppressed terms proportional to the (from below) unconstrained neutrino mass no limit can be derived and terms proportional $`(ϵ_\alpha ^\beta )^2`$ can be neglected. The limits on the remaining non–SM couplings derived in s-wave approximation and evaluated “on axis” are (here and in the following $`90\%C.L.`$): $`ϵ_{V+A}^{V+A}<610^7`$, $`ϵ_{VA}^{V+A}<410^9`$, $`ϵ_{S+P}^{S+P}<910^9`$, $`ϵ_{SP}^{S+P}<910^9`$, $`ϵ_{T_R}^{T_R}<110^9`$, $`ϵ_{T_L}^{T_R}<610^{10}`$. These bounds e.g. exclude the possibility to fake the LSND anomaly (and this way accomodate for all neutrino anomalies with only three neutrinos) via the lepton number violating reaction $`\nu _eu_Ld_Re^+`$ in a model-independent way .
For the short range part the hadronic currents have to be contracted with leptonic currents $`j_\alpha =\overline{e}𝒪_\alpha e^C`$, where $`𝒪_\alpha `$ denotes the operators of defined helicty discussed above. In this case the general Lagrangian is
$`={\displaystyle \frac{G_F^2}{2}}m_P^1\{ϵ_1JJj+ϵ_2J^{\mu \nu }J_{\mu \nu }j+ϵ_3J^\mu J_\mu j+ϵ_4J^\mu J_{\mu \nu }j^\nu +ϵ_5J^\mu Jj_\mu `$
$`+ϵ_6J^\mu J^\nu j_{\mu \nu }+ϵ_7JJ^{\mu \nu }j_{\mu \nu }+ϵ_8J_{\mu \kappa }J^{\nu \kappa }j_\nu ^\mu \},`$ (6)
where indices $`\alpha `$ have been suppressed. Since no fundamental tensors exist in renormalizable theories and since the leptonic tensor current vanishes in the s-wave approximation, the contributions proportional to $`ϵ_4`$, $`ϵ_6`$, $`ϵ_7`$, $`ϵ_8`$ can be neglected. The remaining terms are constrained as follows : $`ϵ_1<310^7`$, $`ϵ_2<210^9`$, $`ϵ_3<410^8/110^8`$ ($`VAVA/VAV\pm A`$), $`ϵ_5<210^7`$.
## 4 Left-right symmetry, R-parity violation, Leptoquarks
In this section we apply the general discussion above to specific theories of physics beyond the SM.
Left–Right–Symmetric Models: In left–right symmetric models the left–handedness of weak interactions is explained as due to the effect of different symmetry breaking scales in the left– and in the right–handed sector. $`0\nu \beta \beta `$ decay proceeds through exchange of the heavy right–handed partner of the ordinary neutrino between right-handed W vertices, leading to a limit of
$$m_{W_R}1.4\left(\frac{m_N}{1TeV}\right)^{(1/4)}TeV.$$
(7)
Including a theoretical limit obtained from considerations of vacuum stability one can deduce an absolute lower limit on the right–handed W mass of
$$m_{W_R}1.4TeV.$$
(8)
Supersymmetry: While in the minimal supersymmetric extension (MSSM) R–parity is assumed to be conserved, there are no theoretical reasons for $`R_p`$ conservation and several GUT and Superstring models require R–parity violation in the low energy regime. In this case $`0\nu \beta \beta `$ decay can occur through Feynman graphs involving the exchange of superpartners as well as $`R_P/`$–couplings $`\lambda ^{^{}}`$ . The half–life limit of the Heidelberg–Moscow experiment leads to bounds in a multidimensional parameter space
$$\lambda _{111}^{^{}}4\times 10^4\left(\frac{m_{\stackrel{~}{q}}}{100GeV}\right)^2\left(\frac{m_{\stackrel{~}{g}}}{100GeV}\right)^{1/2}$$
(9)
(for $`m_{\stackrel{~}{d}_R}=m_{\stackrel{~}{u}_L}`$), which are the sharpest limits on $`R_P/`$–SUSY.
$`0\nu \beta \beta `$ decay is not only sensitive to $`\lambda _{111}^{^{}}`$. Taking into account the fact that the SUSY partners of the left- and right–handed quark states can mix with each other, new diagrams appear in which the neutrino-mediated double beta decay is accompanied by SUSY exchange in the vertices . A calculation of previously neglected tensor contributions to the decay rate allows to derive improved limits on different combinations of $`\lambda ^{^{}}`$ . Assuming the supersymmetric mass parameters of order 100 GeV, the half life limit of the Heidelberg–Moscow Experiment implies: $`\lambda _{113}^{^{}}\lambda _{131}^{^{}}310^8`$, $`\lambda _{112}^{^{}}\lambda _{121}^{^{}}110^6`$
In addition, stringent bounds on coupling products can be derived directly from the effective mass bound eq. 4, since R-parity violating interactions will produce neutrino Majorana masses on loop level. It implies $`\lambda _{133}^{^{}}\lambda _{133}^{^{}}<510^8`$, $`\lambda _{132}^{^{}}\lambda _{123}^{^{}}<110^6`$, $`\lambda _{122}^{^{}}\lambda _{122}^{^{}}<310^5`$, $`\lambda _{133}\lambda _{133}<910^7`$, $`\lambda _{132}\lambda _{123}<210^5`$, $`\lambda _{122}\lambda _{122}<210^4`$.
In the case of R–parity conserving SUSY, based on a theorem proven in , the $`0\nu \beta \beta `$ mass limits can be converted in sneutrino Majorana mass term limits being more restrictive than what could be obtained in inverse neutrinoless double beta decay and single sneutrino production at future linear colliders (NLC) .
Leptoquarks: Leptoquarks are scalar or vector particles coupling both to leptons and quarks, which appear naturally in GUT, extended Technicolor or Compositeness models. The mixing of different multiplets by introducing a leptoquark–Higgs coupling would lead to a contribution to $`0\nu \beta \beta `$ decay . Combined with the half–life limit of the Heidelberg–Moscow experiment bounds on effective couplings can be derived . Assuming only one lepton number violating $`\mathrm{\Delta }L=2`$ LQ–Higgs coupling unequal to zero and the leptoquark masses not too different, one can derive from this limit either a bound on the LQ–Higgs coupling
$`Y_{LQHiggs}=(few)10^6`$ (10)
or a limit excluding leptoquarks with masses in the range of $`𝒪(200GeV)`$. Assuming $`Y_{LQ}𝒪(1)`$ leptoquark masses should be larger than (few) 10 TeV.
## 5 Violations of the equivalence principle and Lorentz invariance
Special relativity and the equivalence principle can be considered as the most basic foundations of the theory of gravity. However, string theories may allow for or even predict the violation of these laws. Such effects in the neutrino sector have been extensively studied in the framework of neutrino oscillations . A typical feature of the violation of Lorentz invariance (VLI) is that different species of matter may have characteristic maximal attainable velocities. The quantity $`\delta v`$ provides an observable for VLI. The corresponding quantity describing violations of the equivalence principle (VEP) is the difference of characteristic couplings $`\delta g`$ to the gravitational potential $`\varphi `$. While previous studies of neutrino oscillations are restricted to the region of large mixing of velocity/gravitational and flavor eigenstates, $`0\nu \beta \beta `$ decay provides a bound in the previously unconstrained region of zero mixing : $`\delta v<410^{16}`$, $`\varphi \delta g<410^{16}`$.
## 6 WIMP Dark Matter Search with Double Beta Experiments
Weakly interacting masssive particles (WIMPs) such as the lightest supersymmetric particle (LSP) are major candidates for the cold component of nonbaryonic dark matter in the universe. Due to its low background properties double beta technology can also find applications in the search for direct detection of WIMPs. The Heidelberg–Moscow Experiment, without being specially designed for this purpose, gave the most stringent limits on WIMPs for several years . New results with 0.69 kg y of measurement reached a background level of 0.042 cts/(kg d keV) in the region between 15 keV and 40 keV. The derived limit excludes WIMPS with masses greater than 13 GeV and cross sections as low as $`1.1210^5`$ pb. These are the most stringent limits on spin-independent interactions using only raw data . The GENIUS experiment would allow to test almost the entire MSSM parameter space already in a first step using only 100 kg of enriched or even natural Ge .
## 7 Summary
Neutrinoless double beta decay and dark matter search belong to the most sensitive approaches with great perspectives to test particle physics beyond the SM.
The possibilities to use $`0\nu \beta \beta `$ decay (and the most sensitive Heidelberg–Moscow experiment) for constraining neutrino masses, new interactions beyond the standard model, and violations of Lorentz invariance and the equivalence principle have been reviewed. Experimental limits on $`0\nu \beta \beta `$ decay are not only complementary to accelerator experiments, neutrino oscillations and cosmological precision measurements, but at least in some cases competitive or superior to the best existing or planned approaches. Direct WIMP detection experiments can compete with recent and future accelerator experiments in the search for SUSY and experiments using double beta technology belong to the most promising approaches in this field of research. A further large breakthrough, both for double beta decay and dark matter search, will be possible realizing the GENIUS proposal, which would improve the obtained limits by up to 1-2 orders of magnitude. |
warning/0002/hep-th0002197.html | ar5iv | text | # Witten-Veneziano Relation for the Schwinger Model
## 1 Euclidean path integral approach on the torus
The topological susceptibility $`\chi _{top}`$ in the theory on the two-dimensional Euclidean torus $`𝒯`$ (with circumferences $`L_1,L_2`$ and volume $`V=L_1L_2`$) is defined as follows
$$\chi _{top}=_𝒯q(x)q(0)d^2x,$$
(3)
where $`q(x)`$ is the topological charge density
$$q(x)=\frac{e}{2\pi }F_{12}(x)\frac{e}{2\pi }E(x)$$
(4)
and
$$q(x)d^2x=\frac{e}{2\pi }E(x)d^2x=k,k\text{}$$
(5)
is the integer-valued topological charge or instanton number in the two-dimensional $`U(1)`$ gauge theory, $`E(x)=_1A_2(x)_2A_1(x)`$ is the field strength. Here and below the spacetime-integrals extend over the torus $`𝒯`$. From Eqs (3) and (4) we obtain
$`\chi _{top}={\displaystyle \frac{e^2}{(2\pi )^2}}{\displaystyle E(x)E(0)d^2x}.`$ (6)
The expectation values $`O(A)`$ of an operator $`O(A)`$ which depends only on the gauge field $`A_\mu (x)`$ are different in the theory with fermions, e.g. the SM, and in the pure electrodynamics on the torus. After integration with respect to fermions in the SM we have
$$O(A)=\frac{1}{Z}_{𝒜_0}𝒟AdetD/_AO(A)e^{\frac{1}{2}{\scriptscriptstyle d^2xE^2(x)}},$$
(7)
where $`D/_A`$ is a Dirac operator in the external electromagnetic field. As a consequence of the index theorem, there are zero modes in the sectors with instanton number $`k0`$ . Therefore the renormalized $`detD/_A`$ is unequal to zero only in the trivial sector $`𝒜_0`$ of the gauge field configurations with vanishing topological charge. $`Z`$ is the partition function in this theory
$$Z=_{𝒜_0}𝒟AdetD/_Ae^{\frac{1}{2}{\scriptscriptstyle d^2xE^2(x)}}.$$
(8)
In pure gauge theory (where expectation values will be denoted by $`\mathrm{}_0`$) all instanton sectors contribute to the expectation value $`O(A)_0`$:
$$O(A)_0=\frac{1}{Z_g}\underset{k\text{}}{}_{𝒜_k}𝒟AO(A)e^{\frac{1}{2}{\scriptscriptstyle d^2xE^2(x)}}$$
(9)
and the partition function reads
$$Z_g=\underset{kZ}{}_{𝒜_k}𝒟Ae^{\frac{1}{2}{\scriptscriptstyle d^2xE^2(x)}}.$$
(10)
In both theories in the sector $`𝒜_k`$ with topological charge $`k`$ the gauge potential has the form
$$A_\mu ^{(k)}(x)=A_\mu ^{(0)}(x)+C_\mu ^{(k)}(x),$$
(11)
where $`A_\mu ^{(0)}(x)`$ is a single valued ”continuous” function on $`𝒯`$ and $`C_\mu ^{(k)}(x)`$ is a global instanton-type potential which in the Lorentz gauge reads
$$C_\mu ^{(k)}(x)=\frac{\pi k}{eV}ϵ_{\mu \nu }x_\nu .$$
(12)
For $`A_\mu ^{(0)}(x)`$ we may use the Hodge decomposition
$$A_\mu ^{(0)}(x)=_\mu a(x)+t_\mu +ϵ_{\mu \nu }_\nu b(x),$$
(13)
where $`_\mu a(x)`$ is pure gauge, $`t_\mu `$ is a (constant) toron field restricted to the dual torus $`0t_\mu T_\mu 2\pi /eL_\mu `$, $`ϵ_{\mu \nu }_\nu b(x)`$ is a curl and $`a(x)`$ and $`b(x)`$ are continuous on $`𝒯`$ and orthogonal to the constant functions: $`a(x)d^2x=b(x)d^2x=0`$ (on the torus the Laplacian $`\mathrm{}_1^2+_2^2`$ is invertible only on functions which integrate to zero).
So the path measure in Eqs.(7) - (10) has a form
$$𝒟A\mathrm{}=𝒟b𝒟a_0^{T_1}𝑑t_1_0^{T_2}𝑑t_2\mathrm{}.$$
(14)
The two-point function $`E(x)E(y)`$ has been calculated in the SM on the torus with the following result:
$$E(x)E(y)=\delta (xy)m^2G_m(xy),$$
(15)
where $`\delta (x)`$ is the $`\delta `$-function on the torus and
$$G_m(x)=\frac{1}{V}\underset{n_1,n_2}{}\frac{e^{2\pi i\left(n_1x_1/L_1+n_2x_2/L_2\right)}}{m^2+\left(\frac{2\pi }{L_1}\right)^2n_1^2+\left(\frac{2\pi }{L_2}\right)^2n_2^2}$$
(16)
is the Greens function of massive scalars on it. From Eqs.(6) and (15) we see that the contact term is
$$P^{(0)}(0)=\frac{e^2}{4\pi ^2}.$$
(17)
As is generally true in gauge theory with massless fermions, the topological susceptibility $`\chi _{top}`$ vanishes in the SM and therefore the relation (2) holds.
Now let us calculate the two-point function $`E(x)E(y)_0`$ in pure electrodynamics. Using the decomposition Eq.(13) we get for the field strength
$$E(x)=\mathrm{}b(x)+\frac{2\pi k}{eV}$$
(18)
and for the action
$$\frac{1}{2}d^2xE^2(x)=\pi \tau k^2+\frac{1}{2}d^2xb(x)\mathrm{}^2b(x),\tau =2\pi /e^2V.$$
(19)
Then
$`E(x)E(y)_0`$ $`=`$ $`\left({\displaystyle \frac{2\pi }{eV}}\right)^2{\displaystyle \frac{k^2e^{\pi \tau k^2}}{e^{\pi \tau k^2}}}+{\displaystyle \frac{𝒟b\mathrm{}b(x)\mathrm{}b(y)e^{\frac{1}{2}{\scriptscriptstyle d^2xb(x)\mathrm{}^2b(x)}}}{𝒟be^{\frac{1}{2}{\scriptscriptstyle d^2xb(x)\mathrm{}^2b(x)}}}}`$ (20)
$`=`$ $`\left({\displaystyle \frac{2\pi }{eV}}\right)^2{\displaystyle \frac{k^2e^{\pi \tau k^2}}{e^{\pi \tau k^2}}}+\delta (xy){\displaystyle \frac{1}{V}},`$
where one sums over all $`k\text{}`$. The presence of the last term in Eq.(20) is due to the fact that $`b(x)`$ does not have a zero mode since it integrates to zero. From Eqs(6) and (20) we get for the topological susceptibility in pure electrodynamics
$`\chi _{top}^{(0)}={\displaystyle \frac{1}{V}}{\displaystyle \frac{k^2e^{\pi \tau k^2}}{e^{\pi \tau k^2}}}.`$ (21)
In pure electrodynamics all instanton sectors contribute to the topological susceptibility. This remains true in infinite volume limit, as we shall see below.
Using the definition of the Jacobi’s $`\theta _3`$ function
$`\theta _3(z|\tau )={\displaystyle \underset{n=\mathrm{}}{\overset{n=\mathrm{}}{}}}e^{\tau \pi n^2+2\pi inz}`$ (22)
we may rewrite Eq.(21) in their terms:
$`\chi _{top}^{(0)}={\displaystyle \frac{1}{4\pi ^2V}}{\displaystyle \frac{\theta _3^{\prime \prime }(0|\tau )}{\theta _3(0|\tau )}}.`$ (23)
## 2 Topological susceptibility in pure electrodynamics on a circle
Pure electrodynamics in two dimensions is defined in a non-trivial way only on a compact space where it has non-trivial gauge invariant solutions. Therefore we consider it again on a circle with circumference $`L_1`$. Manton was the first who considered pure electrodynamics on a circle and showed that in this model there is not a unique canonical quantization, because the representation of the electric field operator contains an arbitrary real parameter $`\theta `$. The Hamiltonian has eigenvalues: $`E_k=\frac{1}{2}L_1e^2\left(k+\theta /2\pi \right)^2,k\text{}`$. In this theory the $`\theta `$ angle (the fractional part of $`\theta /2\pi `$) is a relevant parameter and different values of $`\theta `$ separate different worlds. It corresponds to the famous $`\theta `$ angle of $`SU(n)`$ gauge theories. In contrast, as in the case of quantum chromodynamics with massless fermions , in the SM the angle $`\theta `$ plays non physical role.
The topological susceptibility in this case reads
$`\chi _{top}^{(0)}(0)={\displaystyle \frac{1}{L_1}}{\displaystyle \frac{^2F(\theta )}{\theta ^2}}|_{\theta =0},`$ (24)
where $`F(\theta )=\frac{1}{\beta }\mathrm{log}Z(\theta )`$ is the free energy and $`Z(\theta )`$ the partition function at temperature $`T=1/\beta `$:
$`Z(\theta )={\displaystyle \underset{k}{}}e^{\beta E_k(\theta )}={\displaystyle \underset{k}{}}e^{\frac{1}{2}\beta L_1e^2(k+\theta /2\pi )^2}.`$ (25)
From Eqs.(24) and (25) it follows that
$`\chi _{top}^{(0)}(0)`$ $`=`$ $`{\displaystyle \frac{1}{L_1\beta }}{\displaystyle \frac{Z^{\prime \prime }(0)}{Z(0)}}={\displaystyle \frac{e^2}{4\pi ^2}}\left(1\beta L_1e^2{\displaystyle \frac{k^2e^{\frac{1}{2}\beta L_1e^2k^2}}{e^{\frac{1}{2}\beta L_1e^2k^2}}}\right).`$ (26)
Now we can use the following transformation formula between $`\theta _3`$-functions of zero argument
$`\theta _3(0|\tau )={\displaystyle \frac{1}{\tau }}\theta _3\left(0|1/\tau \right),`$ (27)
and prove that
$`\chi _{top}^{(0)}(0)=\chi _{top}^{(0)},`$ (28)
if we take $`\beta =L_2`$. Thus the path integral approach and canonical approach give the same result for the topological susceptibility.
A systematic comparison between the Hamiltonian approach for the SM on a circle and the Euclidean path integral approach on the torus was done in the forthcoming paper . There it is shown how to obtain Eq.(15) within the Hamiltonian approach.
## 3 The infinite volume limit
In order to consider the limits of infinite volume ($`L_1\mathrm{}`$) and/or the zero temperature $`(L_2\mathrm{})`$ we will use the following expansion :
$`{\displaystyle \frac{\theta _3^{\prime \prime }(0|\tau )}{\theta _3(0|\tau )}}={\displaystyle \frac{2\pi }{\tau }}{\displaystyle \frac{8\pi ^2}{\tau ^2}}e^{\pi /\tau }\left(12e^{\pi /\tau }\right)+\mathrm{}`$ (29)
Then from Eq.(23) we find
$`\chi _{top}^{(0)}={\displaystyle \frac{e^2}{4\pi ^2}}+\mathrm{},`$ (30)
where $`\mathrm{}`$ are terms which disappear if at least one of circumferences $`L_i`$ tends to infinity. Thus we have shown that in this cases the topological susceptibility (23) agrees with the contact term (17).
Acknowledgment: S.A. would like to thank DAAD for a grant which allowed to realize this project. |
warning/0002/hep-th0002191.html | ar5iv | text | # Contents
## 1 Introduction and summary of the results
In the last 4 years the study of black hole solutions in supergravity and string theory has acquired a renewed interest. This is due, essentially, to some successful microscopic entropy countings which have given, for the first time , a statistical interpretation of the Beckestein–Hawking entropy formula in the context of a consistent quantum theory of gravity, such as string theory. The general properties of both supersymmetric and near supersymmetric black holes have been systematically studied and many new results have been obtained.
Despite this progress, it is still not possible nowadays to have a complete control on the common microscopic properties underlying these black holes. Many efforts have been made in order to find a way of making these microscopic entropy countings more based on first principles rather than on case by case computations (see for example ) and to understand which are the actual microscopic degrees of freedom underlying the microscopic/macroscopic matching. However, a definite answer has not been found, yet. In this respect, it would be useful to find a precise description, both at macroscopic and microscopic level, of the so–called generating solution of regular BPS black holes obtained in the framework of superstring (or $`M`$) theory compactified on the torus $`T^6`$ ($`T^7`$).
There are at least three known macroscopic generating solutions in the literature, . However, their structure is too involved to allow a direct construction of the corresponding microscopic configuration. On the other hand, it was proposed in one possible stringy (i.e. microscopic) structure for the generating solution but it was not possible to find its field theory description as a solution of the relevant supergravity theory.
In this paper we fill the above gap and finally find the $`N=8`$ BPS black holes generating solution, giving for it both a macroscopic and a microscopic description in a clear and simple way. As expected, the solution will depend on five independent parameters which will have a precise physical meaning both at macroscopic and microscopic level. Despite the number of independent parameters is five, the number of independent harmonic functions entering the solution will be four. This is somewhat expected and is related to the well known fact that within the five invariants of the $`U`$–duality group of toroidally compactified $`M`$–theory, i.e. $`E_{7(7)}`$, four are moduli–dependent and one is moduli–independent (and proportional to the entropy). Another feature of the generating solution is to have non–trivial axion fields switched on in it (this is related, as we shall see, to the intrinsic property of the central charge eigenvalues evaluated on the solution of being intrinsically complex, i.e. not all imaginary). On the microscopic side this will correspond to have a configuration of D–branes intersecting at non–trivial angles (IIB) or, in a $`T`$–dual picture, a D–brane configuration with non–trivial magnetic fluxes on the D–brane world–volumes (IIA). As expected, the non–trivial four dimensional axions will come from off–diagonal components of the ten dimensional metric $`G_{MN}`$ in the type IIB case and of the antisymmetric tensor $`B_{MN}`$ in the type IIA case, along the world volume of the D–branes (and constant with respect to the world volume coordinates). Switching off the magnetic flux one would recover a four parameter solution of D–branes orthogonally intersecting, which turns out to be a pure dilatonic one.
Once a macroscopic solution it is given, its specific microscopic counterpart it is not uniquely defined: it depends on how the solution is embedded in the original ten dimensional theory. Our choice is to use two suitable R–R embeddings recently defined in , in order to have, as anticipated, a stringy configuration made solely of D–branes. Relying on the structure of the superalgebra central charges and thanks to the geometrical control of the embedding of the solution in the full type II string theory, it will be quite easy to figure out the microscopic systems corresponding to our solution. Moreover, using the tools developed in previous papers ( and especially ) we are able to generate, from the above simple pure D–brane description, any other configuration, even pure NS–NS ones, i.e. those made of solely NS states (as fundamental strings, NS5–branes and KK–states)<sup>2</sup><sup>2</sup>2As shown in , the generating solution of heterotic black holes is also a generating solution for the type II ones. It is the $`U`$–duality group which changes in the two cases and which specifies the $`U`$–duality properties of the solution. In the heterotic case is $`U=SL(2)\times SO(6,22)`$ while in the present case, i.e. type IIA, type IIB or $`M`$–theory compactified on tori ($`T^6`$ or $`T^7`$, respectively), we have $`U=E_{7(7)}`$.. The property for the entropy of being an $`U`$–duality invariant ensures that all these $`U`$–dual configurations share the same entropy. And, as already stressed, the possibility of having a control both at macroscopic and microscopic level of all these configurations, could give some help in unraveling the very conceptual basis of the microscopic entropy counting.
## 2 The macroscopic generating solution
Let us start with the macroscopic description of the generating solution. It has been shown in that the BPS black holes generating solution of toroidally compactified type II string is also solution of a consistent truncation of the relevant $`N=8`$ supergravity effective theory, the so called $`STU`$ model. This is a $`N=2`$ effective model characterized by the graviton multiplet and just three vector multiplets, each of them containing two scalar fields<sup>3</sup><sup>3</sup>3 BPS solutions of $`N=2`$ supergravity have been studied, for instance, in .. Therefore its full bosonic field content is: a graviton, four vector fields (i.e. eight charges, ($`p_\mathrm{\Lambda },q_\mathrm{\Lambda }`$) where $`\mathrm{\Lambda }=0,1,2,3`$) and three complex scalars $`z^i=a_i+\mathrm{i}b_i`$ spanning the Special Kähler manifold $`_{STU}=\left[SL(2,\mathrm{IR})/SO(2)\right]^3`$ ($`b_i`$ being the dilatonic fields and $`a_i`$ the axions). We are not going to describe the detailed structure of the model since it has been described in a complete way in . In fact, we will use the same conventions and notations adopted in those papers. Let us then just briefly summarize the procedure to follow in order to derive the generating solution in the framework of the $`STU`$ model, leaving to the next section the discussion of its embedding in the $`N=8`$ theory. The BPS condition is equivalent to imposing the vanishing of the fermion supersymmetry variation along the Killing spinor $`\xi _a`$ direction:
$`\delta _\xi \text{fermions}`$ $`=`$ $`0`$
$`\gamma ^0\xi _a`$ $`=`$ $`\pm {\displaystyle \frac{Z}{|Z|}}ϵ_{ab}\xi ^b\text{if}a,b=1,2`$ (2.1)
$`Z(z,\overline{z},p,q)`$ being the $`N=2`$ supersymmetry central charge. We adopt the following ansätze for the metric, Killing spinor, scalars and vector field–strengths:
$`ds^2`$ $`=`$ $`e^{2𝒰\left(r\right)}dt^2e^{2𝒰\left(r\right)}d\stackrel{}{x}^2\left(r^2=\stackrel{}{x}^2\right)`$
$`\xi _a(x)`$ $`=`$ $`f(r)ϵ_a`$
$`z^i(x)`$ $`=`$ $`z^i(r)`$
$`F^\mathrm{\Lambda }(r)`$ $`=`$ $`{\displaystyle \frac{p_\mathrm{\Lambda }}{2r^3}}ϵ_{krs}x^kdx^rx^s{\displaystyle \frac{l_\mathrm{\Lambda }(r)}{r^3}}e^{2𝒰(r)}dt\stackrel{}{x}d\stackrel{}{x},\mathrm{\Lambda }=0,1,2,3`$ (2.2)
$`l_\mathrm{\Lambda }`$ being the moduli–dependent electric charges defined in . From the BPS conditions (2.1) we may derive an equivalent system of first order equations for the scalars and metric function $`𝒰`$:
$`{\displaystyle \frac{dz^i}{dr}}`$ $`=`$ $`2\left({\displaystyle \frac{e^{𝒰(r)}}{r^2}}\right)h^{ij^{}}_j^{}|Z(z,\overline{z},p,q)|`$
$`{\displaystyle \frac{d𝒰}{dr}}`$ $`=`$ $`\left({\displaystyle \frac{e^{𝒰(r)}}{r^2}}\right)|Z(z,\overline{z},p,q)|`$ (2.3)
The explicit expression of the right hand side of eq.s (2.3) is quite involved and may be derived from the equations in (computed using the too restrictive condition $`Z=\overline{Z}`$) by multiplying their right hand side by $`\sqrt{Z/\overline{Z}}`$ or by $`\sqrt{\overline{Z}/Z}`$ as far as the equations for the scalars or the one for $`𝒰`$ are concerned, respectively.
In order to characterize the generating solution we need to compute the skew–eigenvalues $`\{Z_\alpha \}=\{Z_{\widehat{i}},Z_4\}`$ ($`\alpha =1,\mathrm{},4`$) of the $`N=8`$ central charge $`Z_{AB}`$ in terms of the supersymmetry and matter central charges $`\{Z,Z_i\}`$ of the $`N=2`$ model, using the following relations:
$`Z=iZ_4,Z^i=h^{ij^{}}_j^{}\overline{Z}=\mathrm{IP}_{\widehat{i}}^iZ^{\widehat{i}}(i,\widehat{i}=1,2,3)`$ (2.4)
where $`\mathrm{IP}_{\widehat{i}}^i=2b_i(r)`$ is the vielbein transforming the rigid indices $`\widehat{i}`$ (the one characterizing the eigenvalues of the $`N=8`$ central charge in its normal form) to the curved indices $`i`$ of the $`STU`$ scalar manifold (see , section 3, for details). Following , the explicit expression of the $`N=8`$ central charge eigenvalues in terms of quantized charges and moduli is given in the appendix. The five $`U`$–duality invariants characterizing a generic solution are the four norms $`|Z_\alpha |`$ and the overall phase $`\mathrm{\Phi }=_\alpha \mathrm{Arg}(Z_\alpha )`$, . These invariants may be combined to give four moduli–dependent invariants and one moduli–independent invariant (the quartic invariant of the $`U`$–duality group). The generating solution is defined as the BPS black hole solution depending on the least number of parameters such that, on the point of the moduli space $`\varphi ^{\mathrm{}}`$ defining the boundary condition at radial infinity of its scalar fields, the five invariants can assume all $`5`$–plets of values (consistent with the positivity condition of the quartic invariant). A necessary condition for the generating solution is thus to depend only on five quantized charges, obtained from the original 8 by suitably fixing the $`SO(2)^3`$ gauge. Our choice for the gauge fixing is $`p_0=q_2=q_3=0`$. As an a–posteriori check that the solution is a generating one, it is necessary to verify that the five invariants computed in the corresponding $`\varphi ^{\mathrm{}}`$ are independent functions of the five remaining charges: $`q_0,q_1,p_1,p_2,p_3`$.
With the above gauge choice the system of first and second order differential equations symplifies considerably and the fixed values for the scalar fields (namely the values the scalars get at the horizon, ) turn out to be the following ones:
$`a_1^{fix}={\displaystyle \frac{q_1p_1}{2p_2p_3}},b_1^{fix}=\sqrt{{\displaystyle \frac{q_0p_1}{p_2p_3}}{\displaystyle \frac{1}{4}}\left({\displaystyle \frac{q_1p_1}{p_2p_3}}\right)^2}`$
$`a_2^{fix}={\displaystyle \frac{q_1}{2p_3}},b_2^{fix}=\sqrt{{\displaystyle \frac{q_0p_2}{p_1p_3}}{\displaystyle \frac{1}{4}}\left({\displaystyle \frac{q_1}{p_3}}\right)^2}`$
$`a_3^{fix}={\displaystyle \frac{q_1}{2p_2}},b_3^{fix}=\sqrt{{\displaystyle \frac{q_0p_3}{p_1p_2}}{\displaystyle \frac{1}{4}}\left({\displaystyle \frac{q_1}{p_2}}\right)^2}`$ (2.5)
Let us now introduce the following harmonic functions:
$`H^i(r)=\mathrm{\hspace{0.17em}1}+{\displaystyle \frac{\sqrt{2}p_i}{r}}\text{with}i=1,2,3`$
$`H_0(r)=\mathrm{\hspace{0.17em}1}+{\displaystyle \frac{\sqrt{2}q_0}{r}}\text{and}H_1(r)=g+{\displaystyle \frac{\sqrt{2}q_1}{r}},g{\displaystyle \frac{q_1}{p_1+p_2}}`$ (2.6)
where the above value for the parameter $`g`$ is fixed by supersymmetry (first order equation for $`𝒰`$).
One can now see that the following ansätze for the $`a_i(r)`$, the $`b_i(r)`$ and the scalar function $`𝒰(r)`$:
$`a_1={\displaystyle \frac{H_1H^1+gH^2}{2H^2H^3}},b_1=\sqrt{{\displaystyle \frac{H_0H^1}{H^2H^3}}{\displaystyle \frac{1}{4}}\left({\displaystyle \frac{H_1H^1gH^2}{H^2H^3}}\right)^2}`$
$`a_2={\displaystyle \frac{H_1H^1gH^2}{2H^1H^3}},b_2=\sqrt{{\displaystyle \frac{H_0H^2}{H^1H^3}}{\displaystyle \frac{1}{4}}\left({\displaystyle \frac{H_1H^1gH^2}{H^1H^3}}\right)^2}`$
$`a_3={\displaystyle \frac{H_1H^1+gH^2}{2H^1H^2}},b_3=\sqrt{{\displaystyle \frac{H_0H^3}{H^1H^2}}{\displaystyle \frac{1}{4}}\left({\displaystyle \frac{H_1H^1gH^2}{H^1H^2}}\right)^2}`$
$`𝒰={\displaystyle \frac{1}{4}}\mathrm{ln}\left(H_0H^1H^2H^3{\displaystyle \frac{1}{4}}(H_1H^1gH^2)^2\right)`$ (2.7)
satisfies both the first and second order differential equations and hence is the solution we were looking for<sup>4</sup><sup>4</sup>4This result is consistent with the analysis in .. The values of the constants characterizing each harmonic function have been chosen in such a way to have 1) asymptotic flat space and 2) asymptotic unitary values for the dilatons $`b_i`$, so to have unitary radii of compactification. The corresponding point in the moduli space at infinity is thus:
$`\varphi ^{\mathrm{}}`$ $``$ $`\{\begin{array}{cc}a_1=a_2=0;a_3=g\hfill & \\ b_i=1\hfill & \end{array}`$ (2.8)
Since the boundary values of the scalar fields at infinity define a bosonic vacuum of the theory, they characterize also the microscopic configuration realizing our solution in the opposite string coupling regime. In next section we shall describe two particularly simple microscopic configurations corresponding to the choice of $`\varphi ^{\mathrm{}}`$ in eq. (2.8).
Notice that the number of truly independent harmonic functions in eq.(2.7) is four, as expected, although the number of independent charges is five: $`q_0,q_1,p_1,p_2,p_3`$. Indeed, from the conditions (2.6), one sees that $`H_1(r)=g\left(H^1(r)+H^2(r)1\right)`$. Finally, according to the ansätze (2.2), the metric has the following form:
$`ds^2`$ $`=`$ $`\left(H_0H^1H^2H^3{\displaystyle \frac{1}{4}}(H_1H^1gH^2)^2\right)^{1/2}dt^2\left(H_0H^1H^2H^3{\displaystyle \frac{1}{4}}(H_1H^1gH^2)^2\right)^{1/2}d\stackrel{}{x}^2`$
and the macroscopic entropy, according to Beckenstein–Hawking formula, reads:
$$S_{macro}=2\pi \sqrt{q_0p_1p_2p_3\frac{1}{4}(q_1p_1)^2}$$
(2.10)
which is the expected expression for the entropy of a generating solution, .
As anticipated, in order to check that the above solution is indeed a five parameters one, one has to work out the expression of the $`N=8`$ central charge skew–eigenvalues on $`\varphi ^{\mathrm{}}`$. From the explicit expressions of the field dependent central charge in the appendix and from (2.8) it follows that:
$`Z_1(\varphi ^{\mathrm{}},p,q)`$ $`=`$ $`{\displaystyle \frac{1}{2\sqrt{2}}}\left[2{\displaystyle \frac{q_1p_1}{p_1+p_2}}+i\left(q_0+p_1p_2p_3\right)\right]`$
$`Z_2(\varphi ^{\mathrm{}},p,q)`$ $`=`$ $`{\displaystyle \frac{1}{2\sqrt{2}}}\left[2{\displaystyle \frac{q_1p_1}{p_1+p_2}}+i\left(q_0p_1+p_2p_3\right)\right]`$
$`Z_3(\varphi ^{\mathrm{}},p,q)`$ $`=`$ $`{\displaystyle \frac{1}{2\sqrt{2}}}\left[0+i\left(q_0p_1p_2+p_3\right)\right]`$
$`Z_4(\varphi ^{\mathrm{}},p,q)`$ $`=`$ $`{\displaystyle \frac{1}{2\sqrt{2}}}\left[0+i\left(q_0+p_1+p_2+p_3\right)\right]`$ (2.11)
From the above equations it is clear that the 5 invariant quantities $`|Z_\alpha (\varphi ^{\mathrm{}},p,q)|`$, $`\mathrm{\Phi }(\varphi ^{\mathrm{}},p,q)`$ are independent functions of the five charges $`q_0,q_1,p_1,p_2,p_3`$.
Switching off the fifth parameter $`q_1`$ our solution becomes exactly the four parameters one studied in (indeed $`H_1(r)|_{q_1=0}=0\text{and}g(q_1=0)=0`$). Indeed putting $`q_1`$ to zero the central charges $`Z_\alpha `$ become pure imaginary and the axion fields vanish uniformly.
## 3 The microscopic description of the generating solution
The microscopic counterpart of a four dimensional macroscopic solution is of course not uniquely defined. Indeed it depends on the interpretation of the four dimensional fields describing the supergravity solution in terms of dimensionally reduced ten dimensional ones. In the light of the analysis put forward in and then completed in <sup>5</sup><sup>5</sup>5Among the main goals of this analysis is the geometrical characterization, using solvable Lie algebra (SLA) techniques, of the scalar and vector fields in the $`N=8`$ theory in terms of type IIA and IIB fields., the microscopic interpretation of our generating solution can be uniquely defined in terms of the embedding of the $`STU`$ model within the original $`N=8`$ theory (in particular of the embedding of the SLA generating the $`STU`$ scalar manifold inside the SLA parametrized by the $`70`$ scalars of the $`N=8`$ theory). In , two main classes of embeddings of the $`STU`$ model were defined (the embeddings within each class being related by $`S,T`$ dualities): one in which the vector fields derive from NS–NS ten dimensional forms and the other in which the vector fields have a R–R origin ( and the scalar fields are NS–NS, see in particular section 2 of that paper). In the latter class two representative embeddings were considered: if we denote by $`x^4,x^5,\mathrm{},x^9`$ the coordinates of $`T^6`$ and by $`x^0,x^1,x^2,x^3`$ the non–compact space–time coordinates, in one embedding, which was characterized from the type IIB point of view, all the three axions of the model come from metric tensor components ($`G_{45},G_{67},G_{89}`$), as opposite to the other, characterized from a type IIA point of view, in which all the axions come from $`B`$ field components ($`B_{45},B_{67},B_{89}`$). In both cases the three dilatons are related to three combinations of the radii of the torus. The two embeddings are related by an operation of $`T`$–duality on the compact directions $`x^5,x^7,x^9`$. As we shall show in the sequel these two embeddings provide an interpretation of the generating solution (2.7) in terms of two $`T`$–dual microscopic configurations: a system of D3–branes at angles (type IIB embedding) and a system of D0 and D4–branes (type IIA embedding) with a magnetic flux in the world volume of the latter (giving therefore extra D2 and D0 charge, ). The magnetic flux (or, equivalently, the non–trivial angle in the dual type IIB configuration) will be the microscopic extra degree of freedom related to the fifth parameter $`q_1`$ characterizing the supergravity solution.
Analyzing the two embeddings from a SLA point of view, one can deduce, in the same way as it was done in for the type IIA case (see in particular eq.s (3.14) and (3.15) of that paper), the subset of the weight basis of the $`\mathrm{𝟓𝟔}`$ of $`E_{7(7)}`$ in terms of which the magnetic $`y^n(\varphi )`$ ($`n=0,\mathrm{},3`$) and electric $`x_n(\varphi )`$ dressed charges<sup>6</sup><sup>6</sup>6 The dressed charges are the physical charges of the interacting theory, which take into account the dressing of the D-brane naked charges provided by the moduli. of the two $`STU`$ model truncations are expressed. According to , the symplectic vector of dressed charges $`(y^n,x_n)`$ is defined in the following way:
$`\left(\begin{array}{c}y^n(\varphi )\\ x_n(\varphi )\end{array}\right)`$ $`=`$ $`\mathrm{IL}^1(\varphi )\left(\begin{array}{c}p^n\\ q_n\end{array}\right)`$ (3.1)
where $`\varphi `$ denotes a point in the scalar manifold and $`\mathrm{IL}(\varphi )`$ is the coset representative of the scalar manifold computed in the same point. As far as the type IIA embedding is concerned, a basis of weights for $`(y^n,x_n)`$ was found in , and from table 3 of the same work the correspondent R–R vectors may be read off:
$`\text{type IIA}:`$
$`(y^n)`$ $``$ $`(A_{\mu 456789},A_{\mu 6789},A_{\mu 4589},A_{\mu 4567})`$
$`(x_n)`$ $``$ $`(A_\mu ,A_{\mu 45},A_{\mu 67},A_{\mu 89})`$ (3.2)
By performing a $`T`$–duality along $`x^5,x^7,x^9`$ according to the geometric recipe given in , we may find the corresponding weights for the type IIB embedding and read from table 3 of the same work their R–R interpretation:
$`\text{type IIB}:`$
$`(y^n)`$ $``$ $`(A_{\mu 468},A_{\mu 568},A_{\mu 478},A_{\mu 469})`$
$`(x_n)`$ $``$ $`(A_{\mu 579},A_{\mu 479},A_{\mu 569},A_{\mu 578})`$ (3.3)
Now let us consider our generating solution and compute the dressed charges on the point of the moduli space $`\varphi ^{\mathrm{}}`$ defined in eq. (2.8). Implementing eq.(3.1), one finds:
$`(y^0,y^1,y^2,y^3)`$ $`=`$ $`(0,p_1,p_2,p_3)`$
$`(x_0,x_1,x_2,x_3)`$ $`=`$ $`(q_0,{\displaystyle \frac{p_1q_1}{p_1+p_2}},{\displaystyle \frac{p_1q_1}{p_1+p_2}},0)`$ (3.4)
From the above expressions we may deduce consistent microscopic configurations corresponding to the generating solution with the chosen boundary condition on the scalar fields at infinity $`\varphi ^{\mathrm{}}`$. From the type IIB viewpoint we may think of a system of D3–branes intersecting at non–trivial angles, but in such a way to preserve 1/8 supersymmetry; this can be achieved if the relative rotation between each couple is a $`SU(3)`$ rotation, . The configuration is depicted in table 1.
Using eq.s (3.3) the first three sets of branes ($`N_0,N_1,N_2`$) may be associated with the charges $`x_0,y^1,y^2`$ respectively, i.e. with charge along the 3–cycles $`(579)`$, $`(568)`$, $`(478)`$, while the fourth set, $`N_3`$, with $`y^3`$, i.e. the charge along $`(469)`$. In fact, due to the non–trivial angle $`\theta `$, the fourth set of $`N_3`$ branes induces D3–brane charge on the cycles $`(579)`$ (contributing to $`x_0`$), $`(479)`$ (represented by $`x_1`$) and $`(569)`$ (represented by $`x_2`$).
As far as the type IIA microscopic interpretation is concerned, we may consider the configuration of D0 and D4–branes obtained by $`T`$–dualizing the type IIB one described above along the directions $`x^5,x^7,x^9`$. The corresponding system may be deduced from eq.s (3.2) and (3.4) and consists of a set of coinciding D0–branes with electric charge $`x_0`$ (the minus sign is required by consistency with the construction in and will be discussed in the sequel) and three sets of coinciding D4–branes along the four–cycles $`(6789)`$, $`(4589)`$ and $`(4567)`$ with magnetic charges $`y^1,y^2,y^3`$. In addition there is a magnetic flux (related to the angle $`\theta `$ in the T–dual type IIB configuration, ) switched on the world volume of the latter brane (i.e. along $`(4567)`$). This flux induces an effective D0 charge (contributing to $`x_0`$) and effective D2 charges along the two–cycles $`(45)`$ and $`(67)`$ (represented by $`x_1`$ and $`x_2`$, respectively). The presence of this flux is also consistent with the fact that the axions in the type IIA embedding are interpreted as coming from the $`B_{MN}`$ tensor in ten dimensions. Indeed, let us briefly recall the general argument relating the presence of a flux on one D4–brane with an effective D2–brane charge (electric in our framework) and a non–trivial $`B_{ij}`$ background field. As well known, $`B`$ field components enter non–trivially in the D$`p`$–brane action via the WZ term:
$$\mu _p_{W_{p+1}}\left(Ce^{}\right)_{p+1}$$
(3.5)
where $``$ is the gauge invariant combination $`=2\pi \alpha ^{}F+\widehat{B}`$ ($`\widehat{B}`$ being the pull–back of the $`B`$ field). Hence, from the supergravity point of view, one would indeed expect new charges representing extra D$`(p2)`$ effective charges at the microscopic level as well as non–trivial bulk $`B`$ field components in the solution (see for instance ). In fact, this is precisely what we get. As shown for instance in , for a suitable choice of the flux (which, albeit giving smaller brane charges via world–volume Chern–Simons coupling, modifies the supersymmetry projections imposed by the D–brane background) the above configuration can preserve $`1/8`$ of the original supersymmetry. Consider the D4–brane configuration described above in the general situation in which the magnetic fluxes are non vanishing on all the three planes $`(45)`$, $`(67)`$, $`(89)`$ (i.e. $`_{45},_{67},_{89}0`$). From eq. (3.5) we may deduce the effective (electric) D2 and D0–brane charges. For instance, the effective D2–brane charge along $`(45)`$ (which is the electric dual object of the D4–brane wrapped on $`(6789)`$) is:
$`\text{\# of }D2\text{ brane (along cycle 45)}={\displaystyle \frac{1}{2\pi }}\left({\displaystyle _{T_{67}}}Tr_{67}+{\displaystyle _{T_{89}}}Tr_{89}\right)`$ (3.6)
which in our conventions is represented by the electric dressed charge $`x_1`$. Similarly, we may compute the effective $`D2`$ charges along the other two–cycles. Notice that the only non–vanishing components of $``$ can only be $`_{45},_{67},_{89}`$ (a $`_{56}0`$ component would imply, for instance, a new $`4`$ dimensional magnetic effective charge out of the four at disposal in the central charge normal gauge), and the three axions come precisely from those components of the $`B`$ field.
On our particular solution the D2–brane charge along $`(89)`$, i.e. $`x_3`$, is zero, while $`(45)`$ and $`(67)`$ charges are opposite one to each other ($`x_1=x_2`$), see (3.4). In fact, considering eq.(3.6) written also for the other 2–cycles, one can easily see that switching on a magnetic flux, as we do, only on the D4–branes lying along $`(4567)`$ and not on the other two bunches of D4–branes, is consistent with having no D2–brane charge along $`(89)`$, i.e. $`x_3=0`$, and opposite D2–brane charge along $`(45)`$ and $`(67)`$, i.e. $`x_1=x_2`$. These charges turn out to be proportional to the fifth parameter $`q_1`$: sending $`q_1`$ to zero the fluxes vanish together with the $`B`$ fields (axions) and we recover the four parameter solution discussed in .
Let us now come to our final goal that is to make the macroscopic/microscopic correspondence precise, namely to give the precise matching between the parameters characterizing the microscopic and the macroscopic configurations, respectively:
$`N_0,N_1,N_2,N_3,\theta q_0,q_1,p_1,p_2,p_3`$
The type IIA and IIB D–brane configurations discussed above as the microscopic counterparts of our generating solution, were suggested in as candidates for the microscopic representation of the 5–parameter solution, whose macroscopic description was then missing. Let us focus for the moment on the type IIA embedding. In order to make contact with this literature, let us use an equivalent representation of the dressed charges, related to the central charges by an $`SO(8)`$<sup>7</sup><sup>7</sup>7The subgroup of $`SU(8)`$ which does not “mix” electric and magnetic charges. transformation:
$$z_{ij}=x_{ij}+iy^{ij}=\frac{1}{\sqrt{2}}\left(\mathrm{\Gamma }^{AB}\right)_{ij}Z_{AB}$$
(3.7)
where the couple $`(ij)`$ indicizes the two times antisymmetric representation of $`SO(8)`$. The real and imaginary parts of $`z_{ij}`$ are the $`N=8`$ electric and magnetic dressed charges in the basis of weights of the $`\mathrm{𝟓𝟔}`$ of $`E_{7(7)}`$ defined in and listed in table 3 of the same paper.
When $`Z_{AB}`$ is skew–diagonal the matrix $`\mathrm{\Gamma }\left(\mathrm{\Gamma }^{AB}\right)_{ij}`$ has the form:
$$\mathrm{\Gamma }=\left(\begin{array}{cccc}\mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.17em}1}& 1& 1& \mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.17em}1}\\ 1& \mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.17em}1}& 1& \mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.17em}1}\\ 1& 1& \mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.17em}1}& \mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.17em}1}\\ 1& 1& 1& 1\end{array}\right)$$
(3.8)
the electric ($`x`$) and magnetic ($`y`$) charges are non vanishing only for $`(ij)`$ equal to $`(12)`$,$`(34)`$,$`(56)`$ and $`(78)`$. From eq.s (2.11) and (3.7) we may read off the values of the charges $`x_{ij}`$ and $`y^{ij}`$:
$`x_{78}+iy^{78}`$ $``$ $`{\displaystyle \frac{1}{\sqrt{2}}}(Z_1+Z_2+Z_3+Z_4)=\mathrm{\hspace{0.33em}0}+iq_0`$
$`x_{12}+iy^{12}`$ $``$ $`{\displaystyle \frac{1}{\sqrt{2}}}(Z_1Z_2Z_3+Z_4)={\displaystyle \frac{q_1p_1}{p_1+p_2}}ip_1`$
$`x_{34}+iy^{34}`$ $``$ $`{\displaystyle \frac{1}{\sqrt{2}}}(Z_1+Z_2Z_3+Z_4)={\displaystyle \frac{q_1p_1}{p_1+p_2}}ip_2`$
$`x_{56}+iy^{56}`$ $``$ $`{\displaystyle \frac{1}{\sqrt{2}}}(Z_1Z_2+Z_3+Z_4)=\mathrm{\hspace{0.33em}0}ip_3`$ (3.9)
The relation between this representation of the dressed charges and the one in eq.(3.4), which is related to the choice of the $`STU`$ model in which the generating solution has been worked out, is the following:
$`\{y^{12},y^{34},y^{56},y^{78}\}`$ $`=`$ $`\{y^1,y^2,y^3,x_0\}`$
$`\{x_{12},x_{34},x_{56},x_{78}\}`$ $`=`$ $`\{x_1,x_2,x_3,y^0\}`$ (3.10)
where symplectic transformation $`x_0=y^{78}`$, $`y^0=x_{78}`$, as discussed in , is related to the feature of our $`STU`$ model (both in type IIA and type IIB cases) of being embedded non–perturbatively in the larger $`N=8`$ theory and therefore is required in order for the truncation to be consistent<sup>8</sup><sup>8</sup>8Indeed the dimensionally reduced R–R vector $`A_\mu `$ in the $`STU`$ model is an electric potential, while it is magnetic in the $`N=8`$ from the type IIA viewpoint, see table 3 of .. The actual D0–brane effective charge is thus $`y^{78}=q_0`$.
The precise correspondence between the dressed charges (in the two representations) and the parameters associated with the microscopic configurations previously discussed (that is those characterizing the type IIB configuration of table 1: $`N_0,N_1,N_2,N_3,\theta `$) is represented in table 2.
Finally, according to relations (3.4) or (3.9) and table 2 we finally get the precise macroscopic/microscopic correspondence<sup>9</sup><sup>9</sup>9Our normalizations are the following. In general the $`4`$ dimensional charge of a wrapped D$`p`$–brane is $`Q_p=\widehat{\mu }_pV_p/\sqrt{V_6}`$ where $`\widehat{\mu }_p=\sqrt{2\pi }(2\pi \sqrt{\alpha ^{}})^{3p}`$ is the normalized D$`p`$–brane charge density in ten dimensions. Provided the asymptotic values of the dilatons $`b_i(r)`$, which parameterize the radii of the compactifying torus and which has been taken to be unitary (see previous section), it turns out that, in units where $`\alpha ^{}=1`$, the four dimensional fundamental quanta of charge for any kind of (wrapped) D$`p`$–brane is equal to $`\sqrt{2\pi }`$ and our quantized charges $`(p_\mathrm{\Lambda },q_\mathrm{\Lambda })`$ have been taken in units of that quanta, i.e. they are integer valued.:
$`N_0=q_0{\displaystyle \frac{(q_1p_1)^2}{p_3(p_1+p_2)^2}},N_1=p_1,N_2=p_2,N_3\mathrm{cos}^2\theta =p_3,\mathrm{tan}\theta ={\displaystyle \frac{q_1p_1}{p_3(p_1+p_2)}}`$
(3.11)
Through equations (3.11) all the microscopic parameters, namely $`N_0,N_1,N_2,N_3`$ and the angle $`\theta `$, are expressed in terms of the quantized charges $`(p_\mathrm{\Lambda },q_\mathrm{\Lambda })`$ characterizing the macroscopic solution and this finally allows us to characterize “quantitatively” its microscopic structure. An alternative microscopic system which could reproduce our generating solution can be given from the $`M`$–theory point of view and consists of three M5–branes with magnetic flux on their world volumes and intersecting on a (compact) line along which $`N_0`$ units of momentum has been put.
With the above definitions, the expression of the $`E_{7(7)}`$ quartic invariant $`J_4`$ in the $`(y^{ij},x_{ij})`$ basis,, is:
$`J_4=4\left(x_{78}x_{12}x_{34}x_{56}+y^{78}y^{12}y^{34}y^{56}\right)\left(x_{78}y^{78}+x_{12}y^{12}+x_{34}y^{34}+x_{56}y^{56}\right)^2`$
$`+\mathrm{\hspace{0.17em}4}\left(x_{78}y^{78}x_{12}y^{12}+x_{78}y^{78}x_{34}y^{34}+x_{78}y^{78}x_{56}y^{56}+x_{12}y^{12}x_{34}y^{34}+x_{12}y^{12}x_{56}y^{56}+x_{34}y^{34}x_{56}y^{56}\right)`$
(3.12)
and consequently, upon use of table 2, one can easily work out the expression of the entropy $`S=\pi \sqrt{J_4}`$ written in terms of the microscopic parameters:
$`S_{micro}=2\pi \sqrt{\mathrm{cos}^2\theta \left[N_0N_1N_2N_3{\displaystyle \frac{1}{4}}\mathrm{sin}^2\theta N_3^2\left(N_1N_2\right)^2\right]}`$ (3.13)
A derivation of the above formula via microscopic counting techniques should be performed extending the analysis of to tori (also the results of could possibly shed some light in this direction). However, we do not try to perform it here. Let us just notice that for $`\theta =0`$ one recovers the usual entropy of the four parameters solution whose derivation via microscopic counting has been carried out, for instance, in .
## 4 Discussion
In the present paper we have worked out the generating solution of four dimensional $`N=8`$ BPS black holes in a form which could be easily described, applying the results of , in terms of pure D–brane configurations upon toroidal compactification of string (or $`M`$) theory. As a result we were able to “pinpoint” the precise correspondence between the microscopic parameters characterizing one of these configurations and the supergravity parameters entering the macroscopic description of the solution.
The relevance of this achievement relies on the possibility on one hand to reconstruct the whole 56–parameter $`U`$–duality orbit of $`N=8`$ BPS black holes, by acting on our solution by means of $`E_{7(7)}`$ transformations, and on the other hand to study in a precise fashion the action of dualities on their corresponding microscopic realizations. Starting from the type IIA configuration described in the previous section and performing a $`T`$–duality transformation on the whole $`T^6`$, one ends up, for instance, with a configuration made of $`N_0`$ D6–branes, 3 bunches of ($`N_1,N_2,N_3`$) D2–branes along the planes ($`45`$),($`67`$),($`89`$) plus effective D4–brane charge. But, more generally, we may also unravel the microscopic properties of pure NS–NS black hole solutions in the same orbit, starting from the corresponding embedding of the $`STU`$ model defined in , or even of mixed NS–NS/R–R solutions. The important point is that having now both a macroscopic and a microscopic description of the generating solution one can follow its trasformation throughout the full $`U`$–duality orbit.
In this respect a challenging problem is to recover the expression of eq.(3.13) from a microscopic entropy counting point of view, performed on the corresponding D–brane configuration. Knowing how to act on it by means of $`U`$–duality can help to shed some light on the actual microscopic degrees of freedom of general BPS black holes, since the generating solution encodes, by definition, all of them. This project is left for future work.
Acknowlodgements
We would like to thank M. Billò, P. di Vecchia, P. Frè, T. Harmark and N. Obers for discussions and R. Russo and C. Scrucca for useful email correspondence. We are also greatful to the organizers of the TMR Torino school on “String Theory and Branes physics” during which this work has been brought to an end. We acknowledge partial support by ECC under contracts ERBFMRX-CT96-0045 and ERBFMRX-CT96-0012.
## Appendix A
The general expression of the $`N=8`$ central charge eigenvalues $`Z_\alpha `$ in terms of the scalar fields and the charges characterizing the $`STU`$ model can be worked out making explicit the first order equations (2.3) taking into account relations (2.4). Following we have the following:
$`ReZ_1`$ $`=`$ $`{\displaystyle \frac{1}{2\sqrt{2b_1b_2b_3}}}(b_1q_1b_2q_2b_3q_3+(a_2a_3b_1a_1a_3b_2a_1a_2b_3b_1b_2b_3)p_0+`$
$`+(a_3b_2+a_2b_3)p_1+(a_1b_3a_3b_1)p_2+(a_1b_2a_2b_1)p_3)`$
$`ImZ_1`$ $`=`$ $`{\displaystyle \frac{1}{2\sqrt{2b_1b_2b_3}}}(a_1q_1+a_2q_2+a_3q_3+(a_1a_2a_3+a_3b_1b_2+a_2b_1b_3a_1b_2b_3)p_0+`$
$`(a_2a_3b_2b_3)p_1(a_1a_3+b_1b_3)p_2(a_1a_2+b_1b_2)p_3+q_0)`$
$`ReZ_2`$ $`=`$ $`(1,2,3)(2,1,3)`$
$`ImZ_2`$ $`=`$ $`(1,2,3)(2,1,3)`$
$`ReZ_3`$ $`=`$ $`(1,2,3)(3,2,1)`$
$`ImZ_3`$ $`=`$ $`(1,2,3)(3,2,1)`$
$`ReZ_4`$ $`=`$ $`{\displaystyle \frac{1}{2\sqrt{2b_1b_2b_3}}}(b_1q_1+b_2q_2+b_3q_3+(a_2a_3b_1+a_1a_3b_2+a_1a_2b_3b_1b_2b_3)p_0+`$
$`(a_3b_2+a_2b_3)p_1(a_3b_1+a_1b_3)p_2(a_2b_1+a_1b_2)p_3)`$
$`ImZ_4`$ $`=`$ $`{\displaystyle \frac{1}{2\sqrt{2b_1b_2b_3}}}(a_1q_1+a_2q_2+a_3q_3+(a_1a_2a_3a_3b_1b_2a_2b_1b_3a_1b_2b_3)p_0+`$ (A.1)
$`(a_2a_3b_2b_3)p_1(a_1a_3b_1b_3)p_2(a_1a_2b_1b_2)p_3+q_0)`$
where it is meant that all axions and dilatons are $`r`$–dependent, i.e. $`a_i=a_i(r)`$ and $`b_i=b_i(r)`$. |
warning/0002/astro-ph0002318.html | ar5iv | text | # ROSAT HRI catalogue of X-ray sources in the LMC regionTable 4 is only available in electronic form at the CDS via anonymous ftp to cdsarc.u-strasbg.fr (130.79.128.5) or via http://cdsweb.u-strasbg.fr/Abstract.html
## 1 Introduction
The Magellanic Clouds (MCs) as the nearest galaxies to the Milky Way allow us to resolve their stellar content in various wavelength bands. X-ray observations combined with optical and radio data can be used to investigate the physical properties of individual X-ray sources as well as the statistical properties of different source classes in a galaxy as a whole. The quantitative and positional distribution of X-ray sources in the MCs will help us to understand the unresolved X-ray emission from more distant galaxies.
After the first observation of X-ray emission from the MCs in 1968 (Mark et al. 1969) four permanent (LMC X-1, X-2, X-3, and X-4, Leong et al. 1971; Giacconi et al. 1972) and few transient X-ray sources were found in the LMC in several satellite missions (UHURU, SAS-3, Copernicus, Ariel-V, HEAO-1). An extensive pointed survey of the LMC was performed by the Einstein Observatory between 1979 and 1981. The two detectors on board this satellite, the Imaging Proportional Counter and the High Resolution Imager, were sensitive enough to detect X-ray binaries, SSSs, and SNRs at the distance of the LMC (55kpc). Long et al. (1981) published a list of 97 discrete X-ray sources in the direction of the LMC and the same data was re-analyzed by Wang et al. (1991) finally giving a list of 105 sources. 54 discrete X-ray sources were identified with objects in the LMC, most of the remaining sources were associated with foreground stars and background AGN. In EXOSAT observations few additional X-ray sources were found (Jones et al. 1985; Pakull et al. 1985; Pietsch et al. 1989).
The next thorough survey of the LMC was made by ROSAT in the energy range of 0.1 – 2.4 keV (Trümper 1982). From 1990 to 1998 ROSAT performed more than 700 pointed observations in a 10 by 10 degree field centered on the LMC. Haberl & Pietsch (1999b, hereafter HP99b) analyzed 212 PSPC observations and created a catalogue of 758 X-ray sources.
In this work results of the analysis of the ROSAT HRI data of the LMC are presented. A description of the HRI detector can be found in David et al. (1996). A source catalogue was obtained in a similar way as in HP99b and many sources were identified by cross-correlating the source list with other existing catalogues. With the help of known properties of different source classes we looked for new candidates for SNRs, stars, and hard X-ray sources which mainly consist of X-ray binaries and absorbed background AGN.
## 2 ROSAT HRI data
### 2.1 Data analysis
The LMC was observed by the ROSAT HRI in more than 500 pointings during the operational phase of ROSAT between 1990 and 1998. 543 observations with exposure times of 50 to 110000 s (Fig. 1) in a field of 10° x 10° around RA = 05<sup>h</sup> 25<sup>m</sup> 00<sup>s</sup>, Dec = –67° 43′ 20″ (J2000.0) were used for the analysis. The analysis was carried out using three detection methods available in EXSAS (Zimmermann et al. 1994). For each pointing X-ray sources were searched using the sliding window methods with local background and with a spline fitted background map. The resulting detection lists were merged and a maximum likelihood algorithm was performed on this list. Sources were accepted if their likelihood of existence was larger than 10.0, i.e. the existence probability was higher than P = 1 – exp(–ML<sub>exi</sub>) = 1 – 4.5 $``$ 10<sup>-5</sup>, and their telescope off-axis angle smaller than 15′ during the observation.
For point and point like sources the source extent was determined by the maximum likelihood technique fitting the source intensity distribution with a Gaussian profile. The count rates resulting from this calculation are correct only for sources with small extent and a brightness profile peaking in the center. For extended sources like SNRs with ringlike structure the net count rates were determined interactively by integrating the counts within a circle around the source. For the background the counts were averaged in a ring around the source distant enough not to be influenced by the source emission.
In order to increase the sensitivity HRI observations with pointing directions within a radius of 1′ were merged after adjusting their position. This was possible for 56 different regions in the LMC. Source detection was also performed on these data and additional faint sources were found which were not detectable in single pointings.
The final source lists obtained for each pointing and co-added observations were merged to one list and multiple detections of a source were reduced to one detection for each source. For this purpose the detection with the smallest positional error was chosen. After screening manually in order to delete spurious detections like knots in extended emission, the catalogue finally contains 397 distinct sources.
### 2.2 Positional corrections and error
ROSAT observations suffer from a systematic positional uncertainty of about 7″ (Kürster 1993). For minimizing this systematic error the coordinates of identified objects were compared to high accuracy positions available in the TYCHO catalogue obtained from the ESA Hipparcos space astrometry satellite (Hoeg et al. 1997) or in the literature. First the X-ray position was corrected to TYCHO coordinates. For sources without any TYCHO counterpart, but identified on the ESO Digitized Sky Survey (DSS) frame with other stars on this frame which were listed in the TYCHO catalogue, more accurate coordinates were calculated for HRI sources by determining the offset between the TYCHO and DSS positions and between the HRI and DSS position. Other sources could be identified with objects in the SIMBAD data base operated at the Centre de Données astronomiques de Strasbourg or in the literature and their positions were corrected after checking their positions on DSS frames. Correction of coordinates for one source implied improved coordinates for all detections of this source in different pointings and for other sources in same pointings. Those secondary corrections again allowed correction of further pointings if the sources were detected several times. Finally for 254 out of 397 sources improved coordinates were determined.
In cases where positional correction was possible the remaining systematic error consists of the error in former optical measurements and the statistical error of the identified source. For not corrected sources the systematic error was set to 7″. The positional error was finally computed as a composite of the statistical uncertainty with 90 % confidence and the systematic error. It is used throughout the paper for the error circle. After the source detection procedure the mean positional error was 8$`\stackrel{}{.}`$3. The coordinate correction reduced the mean positional error of all sources to 6$`\stackrel{}{.}`$4. For position corrected sources the mean positional error is 5$`\stackrel{}{.}`$1.
### 2.3 Correlation with existing catalogues
The catalogue was cross-correlated with the SIMBAD data base and the TYCHO catalogue in order to identify HRI sources. The HRI catalogue contains samples of known SSSs, X-ray binaries, SNRs, Galactic foreground stars, and background AGN. The catalogue was also cross-correlated with the source list from the pointed PSPC observations (HP99b). 138 HRI sources are identical with sources which were detected in PSPC data and thus for most of them the hardness ratios (HR1, HR2) are known. Since the HRI had no spectral resolution no information on the X-ray spectrum could be obtained for HRI sources which are completely new detections. A total of 94 HRI sources were identified with known objects like SSSs, X-ray binaries, SNRs, stars, and background AGN.
With the help of their X-ray properties like extent, extent likelihood, PSPC hardness ratios, X-ray to optical flux ratio (see Sec. 3.2), and X-ray variability 14 previously unknown HRI sources and 7 sources also listed in the PSPC catalogue were newly classified.
The whole source catalogue from HRI observations with the corrected coordinates, final positional error, existence likelihood, HRI count rate, extent, extent likelihood, PSPC count rate and the corresponding PSPC source number with hardness ratios (HP99b) is given in Table 4. For each HRI and PSPC count rate the results for the pointing with the smallest positional error, determined by the maximum likelihood algorithm, were selected. Therefore HRI count rates in the table are representative for one single observation for each source. For extended SNRs the given count rate may correspond to a knot within the source. PSPC count rates are taken from the PSPC catalogue (HP99b) if available. For HRI sources without PSPC detection we derived 2$`\sigma `$ upper limit from the pointing with the highest exposure time. If the source was too close to the rim or the window support structure of the PSPC detector, no count rate is given in Table 4. Neither was it possible to determine PSPC count rates or upper limits for sources located in regions with diffuse emission.
### 2.4 Flux variability
About 80 % of HRI sources were observed more than once and allow time variability studies. For point and point like sources longterm lightcurves were produced with observation-average count rates or upper limits determined by the maximum likelihood algorithm, whereas for extended sources integrated count rates within a circle were used (see Sec. 2.1). For some very bright sources the count rates were integrated in the same way, because an apparent extent resulted from the maximum likelihood algorithm. An apparent extent is computed if the high photon statistics of the bright sources cause a significant deviation from the assumed model for the point spread function.
A $`\chi ^2`$-test for a constant count rate was performed and the factor between the maximum and minimum flux was computed for each lightcurve. Together with the reduced $`\chi ^2`$ this flux factor was used to characterize variability on long time scales of days to years (see also HP99a). For SNRs we expect constant integrated flux, however the flux factor was in the range of 1.0 to 1.8. This may be caused by different off-axis angles and/or different extraction of the extended source. Therefore variations below a factor of 2.0 should be handled with care as they might indicate no real variability but false integration of the source flux because of the extent or existence of a nearby bright source.
In order to obtain a complete lightcurve of the ROSAT observations, also PSPC count rates and upper limits were calculated for the HRI sources. In Figures 2 a – c the PSPC to HRI count rate conversion factor is plotted over N<sub>H</sub> = 10<sup>20</sup> – 10<sup>23</sup> cm<sup>-2</sup> for three different spectral models. SSSs with a soft black body spectrum can be modeled with T = 10.0 – 50.0 eV and galactic N<sub>H</sub> = 10<sup>20</sup> – 10<sup>21</sup> cm<sup>-2</sup> in the direction of the LMC. XBs in general show a power law spectrum with N<sub>H</sub> up to 10<sup>22</sup> cm<sup>-2</sup> because of intrinsic absorption N<sub>H</sub>. So for most of the point and point like X-ray sources PSPC count rates can be converted into HRI count rates by dividing by a typical value of 3, though for very soft sources this scale factor can be larger. Sources in regions with extended emission (e.g. 30 Dor or N44) or close to another source can not always be resolved in PSPC data and may result in false large converting factor.
$`\chi ^2`$ and the flux factor were again calculated for all lightcurves including PSPC count rates (divided by 3.0) and upper limits. Finally 26 sources show significant variability with reduced $`\chi ^2>`$ 5 corresponding to a probability $`>`$ 0.9999 (see Table 1). Four of them are new classified HRI sources (for sources No 49 and 364 see Sec. 3.2.4, for No 300 and 313 see Sec. 3.2.2). As example the lightcurve of source No 49, a new HRI candidate for a variable X-ray binary or AGN is shown in Fig. 3.
PSPC count rates were determined in as many pointings as possible. The mean value was calculated from these count rates and compared to the HRI mean count rates (see Fig. 4). The resulting conversion factor is close to 3.0, only variable sources marked with dots show bigger deviation.
## 3 Source classes
In section 3.1 we discuss HRI sources which were identified either with sources already known from literature or with candidates which were found in former X-ray studies and in PSPC data (HP99b). Section 3.2 deals with new classification of HRI sources based on their X-ray properties.
### 3.1 Source identification
For 97 HRI sources out of 138 which were also detected by the PSPC the HRI observation yielded smaller positional error circles and consequently more accurate source positions compared to the PSPC results. Therefore for several sources likely optical counterparts could be determined which was not possible only with PSPC data.
94 HRI sources were identified with known objects in the LMC, foreground stars, or background objects mainly based on their position (see Sec. 2.3). As they comprise different source types X-ray properties characteristic for each source class could be derived from HRI and PSPC data. Table 3 lists HRI sources with identification.
HP99b have shown that extent and extent likelihood as well as the hardness ratios measured by the PSPC have characteristic values for different source classes and can be used as classification criteria.
In Fig. 5 extent and extent likelihood of the HRI sources are shown. The extent was calculated in the maximum likelihood algorithm and so gives the value resulting from fitting Gaussians. Thus in some cases it may not be the extent of the whole source but only of knots which were found within the extended source. Identified SNRs, marked with open squares, are distributed in the region with large extent and high extent likelihood. Crossed squares indicate known SNR candidates and filled squares sources newly classified as SNR candidates in this work. Point sources have lower extent likelihood unless they were extremely bright like AB Dor (No 180), LMC X-1 (No 311), or RX J0439.8-6809 (No 4) where the deviation of the point spread function from the assumed Gaussian profile becomes significant.
#### 3.1.1 Foreground stars
By cross-correlating the HRI source catalogue with SIMBAD and TYCHO catalogues and using the finding charts presented by Schmidtke et al. (1994, hereafter SCF94), Cowley at al. (1997, hereafter CSM97), and Schmidtke et al. (1999, hereafter SCC99) 39 sources were identified with Galactic foreground stars (Table 3). Most of them could also be identified with the help of UBV photometry results presented by Gochermann et al. (1993) and Grothues et al. (1997). On DSS-images there are point sources as very likely optical counterparts at the positions of these HRI sources within the error circle.
Based on hardness ratios of the PSPC observations two point sources were suggested as foreground star candidates by HP99b (No 189 and 349). They were detected in PSPC images and their hardness ratios are within the range characteristic of stars (HP99b). DSS images show an optical point source within the improved HRI error circle in both cases.
#### 3.1.2 Supernova Remnants
Most SNRs in the LMC are extended X-ray sources which could be resolved by the HRI. They typically show extents of about 5″ – 20″ and high extent likelihood ($`>`$ 10.0). A total of 24 known SNRs were observed by the HRI, four HRI sources are identified with known SNR candidates (No 50, 231, 310, and 315). For both No 231 and 310 the measured hardness ratios are typical for SNRs. No 50 has a harder X-ray spectrum with HR1 = 1.00$`\pm `$0.10 and HR2 = 0.34$`\pm `$0.07.
#### 3.1.3 Supersoft sources
SSSs have very soft X-ray spectra and so far seven SSSs have been discovered in the LMC (HP99b). Two of them were sources of the Einstein LMC survey (Long et al. 1981) and five were found with the help of the ROSAT PSPC. In the HRI pointings five LMC SSSs listed in Table 3 were observed and detected with high existence likelihood.
#### 3.1.4 X-ray binaries
Characteristic for most X-ray binaries is the hard X-ray spectrum and flux variability. In HRI observations nine bright sources could be identified with well known massive X-ray binaries (HMXB). The point source RX J0532.7-6926, here No 238, has been suggested to be a low mass X-ray binary (LMXB) candidate by Haberl & Pietsch (1999a, hereafter HP99a) and was also detected by the HRI. In HP99a a lightcurve with PSPC and HRI measurements is presented and variability is discussed in detail. Between 1990 and 1993 the source showed an exponential intensity decay.
#### 3.1.5 AGN
Nine known background AGN with redshifts between 0.06 and 0.44 (SCF94; CSM97; Crampton et al. 1997) were re-identified in the HRI pointings. Because of its positional coincidence with the radio source PKS 0552-640 and its hardness ratios measured by the PSPC the HRI source No 389 was classified as AGN candidate (No 37 in HP99b). On the DSS frame an optical source with m<sub>B</sub> = 16.3 within the HRI error circle is identified as the most likely optical counterpart.
### 3.2 New classifications
The extensive detection list produced from the HRI pointings towards the LMC allowed us to search for new candidates for different source types. In the course of studying the newly discovered HRI sources the following parameters were of prime importance: count rates, source extent, extent likelihood, flux variability, and counterparts in other wavelengths.
In addition to these X-ray properties we calculated the X-ray to optical flux ratio of HRI sources, for which possible optical counterparts could be found. The flux ratio was computed according to the equation log(f<sub>x</sub>/f<sub>opt</sub>) = log(3 $``$ HRI counts/s $`10^{11}`$) + 0.4 m<sub>B</sub> \+ 5.37 (Maccacaro et al. 1988; HP99b). The relation used for PSPC observations in HP99b was applied here for HRI sources converting the HRI count rates to PSPC count rates by multiplying by the factor of 3 which is typical for hard sources. B magnitudes from the USNO-A1.0 Catalogue produced by the United States Naval Observatory (Monet 1996) were used. For several sources the optical counterpart could not be determined uniquely. In such a case the magnitude of the brightest optical object within the error circle was used resulting in lower limits for log(f<sub>x</sub>/f<sub>opt</sub>). For the SNRs log(f<sub>x</sub>/f<sub>opt</sub>) in general gives no quantitative information, but is an indicator that this source class is bright in X-ray (log(f<sub>x</sub>/f<sub>opt</sub>) $`>`$ –1).
As one can see in Fig. 6 stars in general have negative log(f<sub>x</sub>/f<sub>opt</sub>), for AGN it is around zero, and for SSSs and XBs it is mostly positive in particular when they were observed in their X-ray active phase. Combination of f<sub>x</sub>/f<sub>opt</sub> and the hardness ratios provides a tool to exclude foreground stars.
Newly discovered HRI sources which are suggested as candidates for different source classes in this work can be found in Table 3 and are discussed in the following.
#### 3.2.1 SNR candidates
Investigating the extent five HRI sources (No 197, 284, 288, 307, 338) not classified with the help of PSPC observations are suggested as SNR candidates as their extent is larger than 8″ (see Fig. 5). Since they were not detected by the PSPC because of short exposure times there is no spectral information about these sources which might be crucial for further improvement of the classification.
#### 3.2.2 Sources classified as stellar
For 11 HRI sources probable optical counterparts were found within the error circle which are all bright (m$`{}_{B}{}^{}`$ 12.5), and their log(f<sub>x</sub>/f<sub>opt</sub>) is negative ($`<`$ –2.0). For this reason these sources are classified as stellar objects, and in particular the brightest objects are likely foreground stars. Four sources were also observed by the PSPC (No 90, 135, 217, 313), but as the errors of their hardness ratios are large, no spectral information is given.
The lightcurve of No 300 shows a strong decrease of the X-ray emission with a factor of 10 in 2 years indicating that the HRI observations were performed after an emission maximum. The point source in the optical DSS image at the HRI position is very likely the optical counterpart with a B magnitude of m<sub>B</sub> = 12.4 according to the USNO-A1.0 Catalogue and log(f<sub>x</sub>/f<sub>opt</sub>) = –2.75.
#### 3.2.3 LMC stars as candidates for high mass X-ray binaries
Two X-ray point sources detected by the HRI were identified with known LMC O and B stars (No 328, Sanduleak 1970, m<sub>B</sub> = 18.8 and No 332, Brunet et al. 1975, m<sub>B</sub> = 13.6) because of the positional coincidence. With HRI data no variability investigations could be carried out for these X-ray sources, though there exist many pointings in their direction, because they were both detected only once and in other pointings the upper limits were too high for this purpose. But their identification with optically selected LMC stars allows us to classify them as candidates for high mass X-ray binaries.
#### 3.2.4 Sources with hard X-ray spectrum: Candidates for AGN or X-ray binary
With the help of the hardness ratios and other characteristics measured by the HRI like flux variability or f<sub>x</sub>/f<sub>opt</sub> three HRI sources which were also detected by the PSPC could be classified as candidates either for X-ray binary or for AGN.
The point source No 49 shows significant flux variations, as it is shown in Fig. 3, and has a hard and/or highly absorbed X-ray spectrum (HR1 = 1.00$`\pm `$0.71, HR2 = 0.26$`\pm `$0.16). On the DSS image a likely optical counterpart with a B magnitude of 16.4 (according to the USNO-A1.0 Catalogue) is found. Therefore this source has been classified as a candidate either for an X-ray binary or AGN.
Sources No 230 and 364 are further candidates for X-ray binary or AGN as they have a hard and/or absorbed X-ray spectrum (HR1 = 1.00$`\pm `$0.35, HR2 = 1.00$`\pm `$0.98 and HR1 = 1.00$`\pm `$0.21, HR2 = 1.00$`\pm `$0.60 respectively). Since source No 230 has a small positional error a probable optical counterpart can be found on the DSS image. This counterpart is faint (m<sub>B</sub> = 22.6), and we obtain a high log(f<sub>x</sub>/f<sub>opt</sub>) of 1.56. For source No 364 there is a relatively faint optical source (m<sub>B</sub> = 18.2) within the error circle which might be the counterpart (log(f<sub>x</sub>/f<sub>opt</sub>) = 0.43).
Another nine sources detected by the HRI were identified with sources in the PSPC catalogue (HP99b) showing a hard X-ray spectrum. But from the HRI observations no additional information could be obtained. Thus the HRI sources are simply classified as hard X-ray sources because of the hardness ratios of their PSPC detections.
### 3.3 Source distribution
Due to the high spatial resolution of the HRI many sources could be detected both in the outer regions and in the optical bar region of the LMC. In Fig. 7 HRI sources identified with known objects and known candidates are plotted on a grey scale PSPC image (0.1 – 2.4 keV) of the LMC (from HP99b). The sources are located in different regions of the LMC and show no spatial preferences, it is not only background AGN or foreground stars and candidates which are distributed over the whole LMC region. There are still more than 250 non-identified point sources which are homogeneously distributed in all LMC regions which were covered by ROSAT HRI pointings as it is shown in Fig. 8. In contrast, in PSPC observations not many additional sources could be detected in the regions with strong diffuse emission, because the lower spatial resolution hindered in distinguishing between extended and point like emission (HP99b).
The HRI allows to study the extent of the sources to scales of arcseconds. Therefore SNR candidates could be found not only in regions without surrounding diffuse emission. Four out of five newly suggested SNR candidates are located in regions with diffuse emission between 30 Dor and LMC X-1 (see Fig. 8).
Within and around the optical bar region several new stellar sources and candidates for X-ray binary or AGN were found.
## 4 Summary
The analysis of all 543 ROSAT HRI pointed observations performed between 1990 and 1998 with exposure times higher than 50 sec is presented. Using a maximum likelihood algorithm the source detection resulted in a catalogue of 397 sources which was cross-correlated with the SIMBAD data base and the TYCHO catalogue. Further X-ray properties could be obtained for HRI sources contained in the PSPC catalogue of HP99b.
The high spatial resolution of the HRI enabled the identification of 94 HRI sources with well known objects based on the positional coincidence and considering their extent and hardness ratios. The coordinates of most of the identified sources could be improved to more accurate positions and allowed the positional correction of other HRI sources. Thus for 254 sources the systematic error for their coordinates could be reduced to values smaller than 7″ which is the standard systematic error of ROSAT observations.
For different source classes like SSS, X-ray binary, SNR, Galactic stars, and background AGN classification criteria could be derived from the extent and hardness ratios of the identified sources. We looked for flux variability of the sources and for likely optical counterparts. Five newly detected HRI sources were classified as candidates for SNRs because of their extent, two HRI sources which were identified with an LMC O and a B star as HMXB candidates. Eleven sources with probable bright optical counterpart and small X-ray to optical flux ratio are classified as stellar sources. Three sources with hard and/or absorbed X-ray spectrum indicated by the PSPC hardness ratios are likely candidates for X-ray binaries or AGN. Two of the hard X-ray sources show flux variability and for each of these an optical counterpart was found.
With the help of HRI observations many new X-ray sources were found. Further follow-up observations in X-ray, optical, or radio wavelengths with spectral information are needed to characterize these sources in more detail.
###### Acknowledgements.
The ROSAT project is supported by the German Bundesministerium für Bildung und Forschung (BMBF) and the Max-Planck-Gesellschaft. This research has been carried out by making extensive use of the SIMBAD data base operated at CDS, Strasbourg, France. |
warning/0002/hep-th0002211.html | ar5iv | text | # Contents
## 1 Introduction
The spectrum of open strings living on a D-brane anti-D-brane pair of type IIA or IIB string theory contains a pair of tachyonic modes from the Neveu-Schwarz (NS) sector, indicating that the system is unstable . There are general arguments which indicate that the tachyonic potential has a minimum, and this minimum represents the usual vacuum of the closed string theory without any D-brane.<sup>1</sup><sup>1</sup>1Some of these arguments used earlier gauge theory analysis of brane-antibrane systems, see . For this to be true, the negative energy density contribution from the tachyon potential at the minimum must exactly cancel the sum of the tensions of the brane-antibrane pair.
During the last few years it has been realised that type IIA (IIB) string theory contains unstable non-BPS D$`p`$-branes for odd (even) $`p`$ . Any of these D-branes has a tachyonic mode, indicating that the brane is unstable. A straightforward extension of the general arguments for the brane-antibrane system can be used to argue that the tachyonic potential has a minimum, and this minimum represents the usual vacuum of the closed string theory without any D-brane. For this to be true, the negative energy density contribution from the tachyon potential at this minimum must exactly cancel the tension of the non-BPS D-brane.
A similar conjecture exists also for D-branes of bosonic string theory . On any bosonic string D-brane there is a tachyonic open string mode. Indirect arguments, similar to those for the brane-antibrane pair of type II string theories, indicate that the tachyon potential has an extremum whose negative energy density contribution cancels the tension of the D-brane, so that this particular extremum represents the vacuum of closed bosonic string theory without any D-branes. In refs. this phenomenon was studied directly in open bosonic string field theory , following earlier work of ref.. Using the level truncation scheme of ref., ref. showed that including scalars up to level four, the value of the potential at the extremum cancels almost 99% of the D-brane tension. This is a strong indication that this extremum indeed represents vacuum without D-branes. This remarkable cancellation has now been verified to an accuracy of 99.9% by including scalars up to level ten . Evidence to the validity of the level expansion was recently given in . Further evidence for the identification of the tachyonic vacuum has been found in ref. who considered tachyonic lump solutions of the string equations of motion.
In a recent paper , the zeroth order contribution to the tachyon potential on a non-BPS D-brane of type II string theory was computed using the open string field theory action formulated in refs., and was found to contain a minimum at which the potential cancels 60% of the D-brane tension. Unlike the cubic action proposed in , the Wess-Zumino-Witten-like action used in ref. has no problems with contact term divergences .<sup>2</sup><sup>2</sup>2An early attempt at generalizing the analysis of ref. to superstring theory was made in ref.. Although it is not known how to include Ramond (R) sector states into this action in a manifestly SO(9,1) covariant manner, this is not a problem here since the phenomenon of tachyon condensation involves NS sector states only. Of course, it involves the full unprojected NS sector, namely both GSO$`(+)`$ and GSO$`()`$ states. While there is a superspace version of the Wess-Zumino-Witten-like action which is manifestly SO(3,1) super-Poincaré invariant and includes all GSO$`(+)`$ states in the NS and R sectors, it is not yet known how to incorporate GSO$`()`$ states into it.
In this paper, we compute the first-order correction to the tachyon potential on a non-BPS D-brane of type II string theory using the same action as , and find a minimum of the potential at which 85% of the D-brane tension is cancelled by the potential. This provides strong evidence that the approximation scheme is converging as in the bosonic string computation and that the exact tachyon potential has a minimum where the D-brane tension is exactly cancelled. Alternatively, this result can also be viewed as a successful test of the correctness of this superstring field theory action.
Although we carry out the explicit analysis for the non-BPS D-brane, the result also holds for the brane-antibrane system. Indeed, the tachyon potential on the non-BPS D-brane of type IIA (IIB) string theory can be obtained from the tachyon potential on a brane-antibrane system of type IIB (IIA) string theory after restricting the field configuration to a $`Z_2`$ invariant subspace; so the existence of a minimum of the tachyon potential for a non-BPS D-brane corresponding to the vacuum without D-branes also establishes the corresponding result for the brane-antibrane system. This can also be made self-evident by comparing the structure of the string field theory action on the non-BPS D-brane and that on a brane-antibrane pair, both of which we write down explicitly.
Since the tachyon potential on a non-BPS D-brane is invariant under a change of sign of the tachyon field, there are doubly degenerate minima of the potential, and we can construct a kink solution interpolating between these two minima. It has been conjectured that this represents a BPS brane of one lower dimension following a similar conjecture for the brane-antibrane system. We compute numerically the energy density of the kink solution using the tachyon potential, but ignoring string field theory corrections to the tachyon kinetic term, in the same spirit as in a recent paper for bosonic string field theory . The result is 1.03 times the expected answer. Although such a close agreement is likely to be accidental, it is encouraging to note that the mass of the kink even in this crude approximation has the correct order of magnitude. We should also note that the effect of the non-zero tachyon background should be to reduce the kinetic term, since at the minimum of the potential the kinetic term is expected to vanish, so that we have no physical excitations. Thus we expect that once we take into account corrections to the kinetic term, the energy of the kink should be lowered. In this context, it is also encouraging to note that in the analysis of ref. the mass of the lump decreased after taking into account corrections to the tachyon kinetic energy term.
The paper is organised as follows. In section 2 of this paper, we shall review this WZW-like action, discuss in detail its basic ingredients, its gauge invariance, and its application to describe the non-BPS D-brane as well as the brane anti-brane system. In section 3, we shall use the action to compute the zeroth and first order contributions to the tachyon potential and show that the potential has a minimum at 85% of the D-brane tension. In section 4 we discuss the tachyonic kink solution and calculate its mass. We offer some perspective on our results and discuss open questions in section 5. Important details have been provided in the appendices. Appendix A establishes the cyclicity properties of the amplitudes appearing in the string action– this cyclicity is essential for gauge invariance. Appendix B explains the twist properties of the amplitudes– such properties allows us to restrict the multiscalar tachyon field (the space $`_1`$ defined in section 3) to the twist even sector. Appendix C gives a self-contained derivation of the mass of the D-brane described by a string field theory action. Finally, in appendix D we provide details on the computation of the tachyon potential.
## 2 Open superstring field theory
In this section we shall explain and analyze the superstring field theory that describes the dynamics of a non-BPS D-brane of type II string theory. As it will be clear, this string field theory is readily modified to discuss the D-brane anti-D-brane system in superstring theory. In fact, the same calculations give the tachyon potential for both physical situations. As in the case of refs., we shall not restrict our analysis to any specific background, but will assume, for convenience, that all the directions tangential to the D-brane are compact, so that the system has a finite mass.
We will begin by discussing the GSO projected, or GSO$`(+)`$ sector of the open superstring theory formulated in refs., $``$ this would describe the dynamics of NS sector open strings living on a single BPS D-brane. Here the basic structure of the theory will be elaborated. Then we turn to the non-BPS D-brane whose formulation requires incorporating both the GSO$`()`$ sector and the GSO$`(+)`$ sector of the theory. This can be done by attaching internal Chan-Paton matrices to the GSO plus and minus sectors in such a way that the complete string field and the relevant operators satisfy the basic structure of the original GSO$`(+)`$ theory. This device was used in for the analysis of the non-BPS D-brane. Finally, in the last subsection we show how, in addition to the internal Chan-Paton matrices, external Chan-Paton matrices must be tensored to describe the brane-antibrane system.
### 2.1 Superstring field theory on a BPS D-brane
In the formalism of refs., a general off-shell string field configuration in the GSO$`(+)`$ NS sector corresponds to a Grassmann even open string vertex operator $`\mathrm{\Phi }`$ of ghost number 0 and picture number 0 in the combined conformal field theory of a $`c=15`$ superconformal matter system, and the $`b,c,\beta ,\gamma `$ ghost system with $`c=15`$. In terms of the bosonized ghost fields $`\xi ,\eta ,\varphi `$ related to $`\beta ,\gamma `$ through the relations
$$\beta =\xi e^\varphi ,\gamma =\eta e^\varphi ,$$
(2.1)
the ghost number ($`n_g`$) and the picture number ($`n_p`$) assignments are as follows:
$`b:n_g=1,n_p=0,c:n_g=1,n_p=0,`$
$`e^{q\varphi }:n_g=0,n_p=q,`$
$`\xi :n_g=1,n_p=1,\eta :n_g=1,n_p=1.`$
The SL(2,R) invariant vacuum carries zero ghost and picture number. Note that this definition of ghost number agrees with the definition of for states with zero picture, but unlike the definition of , it allows the spacetime-supersymmetry generators to carry zero ghost number. One notable difference from other formulations of open string field theory is that here the string field correspond to vertex operators in the ‘large Hilbert space’ containing the zero mode of the field $`\xi `$.
We shall denote by $`_iA_i`$ the correlation function of a set of vertex operators in the combined matter-ghost conformal field theory on the unit disk with open string vertex operators inserted on the boundary of the disk, without including trace over CP factors. These correlation functions are to be computed with the normalization
$$\xi (z)cc^2c(w)e^{2\varphi (y)}=2.$$
(2.3)
Throughout this paper we shall be working in units where $`\alpha ^{}=1`$. The nilpotent BRST operator of this theory is given by
$$Q_B=𝑑zj_B(z)=𝑑z\left\{c(T_m+T_{\xi \eta }+T_\varphi )+ccb+\eta e^\varphi G_m\eta \eta e^{2\varphi }b\right\},$$
(2.4)
where
$$T_{\xi \eta }=\xi \eta ,T_\varphi =\frac{1}{2}\varphi \varphi ^2\varphi ,$$
(2.5)
$`T_m`$ is the matter stress tensor and $`G_m`$ is the matter superconformal generator. $`G_m`$ is a dimension $`3/2`$ primary field and satisfies:
$$G_m(z)G_m(w)\frac{10}{(zw)^3}+\frac{2T_m}{(zw)}.$$
(2.6)
The normalization of $`\varphi `$, $`\xi `$, $`\eta `$, $`b`$ and $`c`$ are as follows:
$$\xi (z)\eta (w)\frac{1}{zw},b(z)c(w)\frac{1}{zw},\varphi (z)\varphi (w)\frac{1}{(zw)^2}.$$
(2.7)
We denote by $`\eta _0=𝑑z\eta (z)`$ the zero mode of the field $`\eta `$ acting on the Hilbert space of matter ghost CFT.
The string field theory action is given by
$$S=\frac{1}{2g^2}(e^\mathrm{\Phi }Q_Be^\mathrm{\Phi })(e^\mathrm{\Phi }\eta _0e^\mathrm{\Phi })_0^1𝑑t(e^{t\mathrm{\Phi }}_te^{t\mathrm{\Phi }})\{(e^{t\mathrm{\Phi }}Q_Be^{t\mathrm{\Phi }}),(e^{t\mathrm{\Phi }}\eta _0e^{t\mathrm{\Phi }})\},$$
(2.8)
where $`\{A,B\}AB+BA`$, and $`e^{t\mathrm{\Phi }}_te^{t\mathrm{\Phi }}=\mathrm{\Phi }`$ but has been written this way for convenience. This action is defined by expanding all exponentials in formal Taylor series carefully preserving the order of all operators and letting $`\mathrm{}`$ of an ordered sequence of arbitrary vertex operators $`A_1,\mathrm{}A_n`$ be defined as:
$$A_1\mathrm{}A_n=f_1^{(n)}A_1(0)\mathrm{}f_n^{(n)}A_n(0).$$
(2.9)
Here, $`fA`$ for any function $`f(z)`$, denotes the conformal transform of $`A`$ by $`f`$, and
$$f_k^{(n)}(z)=e^{\frac{2\pi i(k1)}{n}}\left(\frac{1+iz}{1iz}\right)^{2/n}\text{for}n1.$$
(2.10)
In particular if $`\phi `$ denotes a primary field of weight $`h`$, then
$$f\phi (0)=(f^{}(0))^h\phi (f(0)).$$
(2.11)
$`Q_B`$ (or $`\eta _0`$) acting on a set of vertex operators inside $``$ corresponds to a contour integral of $`j_B`$ (or $`\eta `$) around the insertion points of these vertex operators on the right hand side of eq.(2.9).
Since we have, in general, non-integer weight vertex operators, we should be more careful in defining $`fA`$ for such vertex operators. Noting that
$$f_k^{(N)}(0)=\frac{4i}{N}e^{2\pi i\frac{k1}{N}}\frac{4}{N}e^{2\pi i(\frac{k1}{N}+\frac{1}{4})},$$
(2.12)
we adopt the following definition of $`f_k^{(N)}\phi (0)`$ for a primary vertex operator $`\phi (x)`$ of conformal weight $`h`$:
$$f_k^{(N)}\phi (0)=\left|\left(\frac{4}{N}\right)^h\right|e^{2\pi ih(\frac{k1}{N}+\frac{1}{4})}\phi (f_k^{(N)}(0)).$$
(2.13)
Since all secondary vertex operators can be obtained as products of derivatives of primary vertex operators, this uniquely defines $`f_k^{(N)}A(0)`$ for all vertex operators.
The geometry of the interaction described in (2.10) is simple. The function $`f_1^{(n)}`$ maps the upper half disk $`|z|1,\mathrm{}(z)>0`$ into the wedge $`|\text{Arg}(f_1^{(n)})|\pi /n,|f_1^{(n)}|1`$, with the puncture $`z=0`$ ending at $`f_1^{(n)}=1`$. With $`k=1,\mathrm{}n`$, we end up gluing $`n`$ such wedges together to form a full unit disk where the $`n`$ vertex operators are inserted at equally spaced points on the boundary. By a further SL(2,C) transformation $`F`$ (e.g. $`F(w)=i(1w)/(1+w)`$) we can map the interior of the unit disk onto the upper half plane. We could use the functions $`g_k^{(n)}(z)=F(f_k^{(n)}(z))`$ instead of $`f_k^{(n)}(z)`$ to define the string field theory action. $``$ will now denote the correlation function of the conformal field theory on the upper half plane, with open string vertex operators inserted on the real axis. As will be shown in appendix D, by a convenient choice of the SL(2,C) transformation $`F`$ we can ensure that $`g_1^{(n)}(0),\mathrm{}g_n^{(n)}(0)`$ are ordered from left to right on the real axis. Also one finds that $`g_k^{(n)}(0)`$ is real and positive for all $`k`$. The prescription (2.13) then corresponds to choosing real, positive values of $`(g_k^{(n)}(0))^h`$ in the expression for the conformal transform of a field $`\mathrm{\Phi }`$ of weight $`h`$. As a double-check of our computations, we shall compute the tachyon potential using both the disk and the UHP prescriptions and compare answers.
The correlator $``$ defined in eq.(2.9) satisfies cyclicity properties. Let $`\mathrm{\Phi }`$ denote any component of the string field, and $`A_1,\mathrm{}A_{n1}`$ denote arbitrary vertex operators (i.e. arbitrary grassmanality, ghost number, etc.). Then
$`A_1\mathrm{}A_{n1}\mathrm{\Phi }=\mathrm{\Phi }A_1\mathrm{}A_{n1},`$
$`A_1\mathrm{}A_{n1}(Q_B\mathrm{\Phi })=(Q_B\mathrm{\Phi })A_1\mathrm{}A_{n1},`$
$`A_1\mathrm{}A_{n1}(\eta _0\mathrm{\Phi })=(\eta _0\mathrm{\Phi })A_1\mathrm{}A_{n1}.`$ (2.14)
The proof of these relations has been given in appendix A.
Note that in this notation the BPZ inner product is given by:
$$A|B=AB,$$
(2.15)
which uses the two punctured disk (eqn.(2.10) with $`n=2`$). We now define the multilinear products $`|A_1A_2\mathrm{}A_n`$ of $`n`$ vertex operators $`A_1,A_2,\mathrm{}A_n`$ through the relation:
$$B|A_1\mathrm{}A_n=BA_1\mathrm{}A_n,$$
(2.16)
for any state $`B|`$. The product $`|A_1A_2`$, computed with (2.10) and $`n=3`$, is simply the associative (non-commutative) star product $`|A_1A_2`$ of . It follows from the geometry of the interaction that the higher products are equivalent to iterated multiplication using the star product: namely, $`|A_1A_2\mathrm{}A_n=|A_1A_2\mathrm{}A_n`$. While the order of the sequence of operators must be preserved, the multiplications in this ket can be done in any order, thanks to the associativity of the star product. It follows that all products associate. From now on we shall denote the product of a set of vertex operators $`A_1,A_2,\mathrm{}A_n`$ by $`A_1A_2\mathrm{}A_n`$.
It will now be shown that this action is invariant under the gauge transformation
$$\delta e^\mathrm{\Phi }=(Q_B\mathrm{\Omega })e^\mathrm{\Phi }+e^\mathrm{\Phi }(\eta _0\mathrm{\Omega }^{}),$$
(2.17)
where the gauge transformation parameters $`\mathrm{\Omega }`$ and $`\mathrm{\Omega }^{}`$ are Grassmann odd, GSO$`(+)`$ vertex operators with $`(n_g,n_p)`$ values $`(1,0)`$ and $`(1,1)`$ respectively. The proof will use the cyclicity relations (LABEL:e1aa) and the following identities:
$$\{Q_B,\eta _0\}=0,Q_B^2=\eta _0^2=0,$$
(2.18)
$$Q_B(\mathrm{\Phi }_1\mathrm{\Phi }_2)=(Q_B\mathrm{\Phi }_1)\mathrm{\Phi }_2+\mathrm{\Phi }_1(Q_B\mathrm{\Phi }_2),\eta _0(\mathrm{\Phi }_1\mathrm{\Phi }_2)=(\eta _0\mathrm{\Phi }_1)\mathrm{\Phi }_2+\mathrm{\Phi }_1(\eta _0\mathrm{\Phi }_2),$$
(2.19)
$$Q_B(\mathrm{})=\eta _0(\mathrm{})=0.$$
(2.20)
Note that in the identities of the second line there are no minus signs necessary as $`Q_B`$ or $`\eta _0`$ go through the string field because the string field is Grassmann even. The identities in the last last line hold because $`Q_B`$ and $`\eta _0`$ are integrals of dimension one currents which can be “pulled” off the boundary and collapsed inside the disk.
Defining $`G=e^\mathrm{\Phi }`$ and using the above identities, one finds that under an arbitrary variation $`\delta G`$,
$$\delta S=\frac{1}{g^2}G^1\delta G\eta _0(G^1Q_BG)$$
(2.21)
where the first term of $`S`$ contributes $`(2g^2)^1G^1\delta G[\eta _0(G^1Q_BG)Q_B(G^1\eta _0G)]`$ to the variation and the second term of $`S`$ contributes $`(2g^2)^1G^1\delta G[\eta _0(G^1Q_BG)+Q_B(G^1\eta _0G)]`$ to the variation. Using the fact that $`S`$ goes to $`S`$ after switching $`\eta _0`$ with $`Q_B`$ and $`G`$ with $`G^1`$, (2.21) can also be written as
$$\delta S=\frac{1}{g^2}G\delta G^1Q_B(G\eta _0G^1)$$
(2.22)
To prove gauge invariance under $`\delta G=G\eta _0\mathrm{\Omega }^{}`$, use (2.21) and pull $`\eta _0`$ off the $`(G^1Q_BG)`$ term. Since $`\eta _0(G^1\delta G)=0`$, $`\delta S=0`$. To prove gauge invariance under $`\delta G=(Q_B\mathrm{\Omega })G`$, use (2.22) and pull $`Q_B`$ off of the $`(GQ_BG^1)`$ term. Since $`Q_B(G\delta G^1)=Q_B(\delta GG^1)=0`$, $`\delta S=0`$. So we have proven invariance of the action under the gauge transformations of (2.17).
The equation of motion for the action is easily derived from (2.21) to be
$$\eta _0(e^\mathrm{\Phi }Q_Be^\mathrm{\Phi })=0.$$
(2.23)
As stated earlier, the string field in the present theory corresponds to vertex operators in the ‘large Hilbert space’ which includes the zero mode of $`\xi `$. However, using the $`\mathrm{\Omega }^{}`$ gauge invariance, we can choose the gauge $`\xi _0\mathrm{\Phi }=0`$. In that gauge, the string field configuration $`\mathrm{\Phi }`$ is in one to one correpondence with vertex operators in the ‘small Hilbert space’ which does not include the zero mode of $`\xi `$. This will be discussed in some detail in section 3.
### 2.2 Superstring field theory on a Non-BPS D-brane
The open string states living on a single non-BPS D-brane are divided into two classes, GSO$`(+)`$ states and GSO$`()`$ states. Since the GSO$`()`$ states are Grassmann odd they cannot be incorporated directly into a string field preserving the algebraic structure reviewed in the previous subsection. This structure can be recovered by tensoring $`2\times 2`$ matrices carrying internal Chan-Paton (CP) indices. These are added both to the vertex operators and to $`Q_B`$ and $`\eta _0`$.
We attach the $`2\times 2`$ identity matrix $`I`$ on the usual GSO$`(+)`$ sector (recall that the Neveu-Schwarz (NS) sector ground state is odd under the projection operator $`(1)^F`$) and the Pauli matrix $`\sigma _1`$ to the GSO$`()`$ sector. The complete string field is thus written as
$$\widehat{\mathrm{\Phi }}=\mathrm{\Phi }_+I+\mathrm{\Phi }_{}\sigma _1,$$
(2.24)
where the subscripts denote the $`()^F`$ eigenvalue of the vertex operator. In addition, we define:
$$\widehat{Q}_B=Q_B\sigma _3,\widehat{\eta }_0=\eta _0\sigma _3.$$
(2.25)
Note that this definition shows that these matrices do not really carry conventional CP indices; had it been so, both $`Q_B`$ and $`\eta _0`$ should have been tensored with the identity matrix, as such operators should not change the sector the strings live in . Finally, we define
$$\widehat{A}_1\mathrm{}\widehat{A}_n=Trf_1^{(n)}\widehat{A}_1(0)\mathrm{}f_n^{(n)}\widehat{A}_n(0),$$
(2.26)
where the trace is over the internal CP matrices. We shall adopt the convention that fields or operators with internal CP factors included are denoted by symbols with a hat on them, and fields or operators without internal CP factors included are denoted by symbols without a hat, as in the previous subsection.
Indeed, with these definitions the cyclicity relations (LABEL:e1aa) given in the previous section now hold as
$`\widehat{A}_1\mathrm{}\widehat{A}_{n1}\widehat{\mathrm{\Phi }}=\widehat{\mathrm{\Phi }}\widehat{A}_1\mathrm{}\widehat{A}_{n1},`$
$`\widehat{A}_1\mathrm{}\widehat{A}_{n1}(\widehat{Q}_B\widehat{\mathrm{\Phi }})=(\widehat{Q}_B\widehat{\mathrm{\Phi }})\widehat{A}_1\mathrm{}\widehat{A}_{n1},`$
$`\widehat{A}_1\mathrm{}\widehat{A}_{n1}(\widehat{\eta }_0\widehat{\mathrm{\Phi }})=(\widehat{\eta }_0\widehat{\mathrm{\Phi }})\widehat{A}_1\mathrm{}\widehat{A}_{n1},`$ (2.27)
where $`\widehat{\mathrm{\Phi }}`$ denotes any component of the the string field, and $`\widehat{A}_1,\mathrm{}\widehat{A}_{n1}`$ denote arbitrary vertex operators. The proof of these relations, as well as those for the unhatted case have been given in appendix A. In addition, we have the analogs of (2.18) holding
$$\{\widehat{Q}_B,\widehat{\eta }_0\}=0,\widehat{Q}_{B}^{}{}_{}{}^{2}=\widehat{\eta }_{0}^{}{}_{}{}^{2}=0,$$
(2.28)
$$\widehat{Q}_B(\widehat{\mathrm{\Phi }}_1\widehat{\mathrm{\Phi }}_2)=(\widehat{Q}_B\widehat{\mathrm{\Phi }}_1)\widehat{\mathrm{\Phi }}_2+\widehat{\mathrm{\Phi }}_1(\widehat{Q}_B\widehat{\mathrm{\Phi }}_2),\widehat{\eta }_0(\widehat{\mathrm{\Phi }}_1\widehat{\mathrm{\Phi }}_2)=(\widehat{\eta }_0\widehat{\mathrm{\Phi }}_1)\widehat{\mathrm{\Phi }}_2+\widehat{\mathrm{\Phi }}_1(\widehat{\eta }_0\widehat{\mathrm{\Phi }}_2),$$
(2.29)
$$\widehat{Q}_B(\mathrm{})=\widehat{\eta }_0(\mathrm{})=0.$$
(2.30)
The reason no extra signs appear in the middle equation is clear, when the string field is Grassmann odd the sign arising by moving $`Q_B`$ across the vertex operator is cancelled by having to move $`\sigma _3`$ across $`\sigma _1`$.
Given that the relations satified by the hatted objects are the same as those of the unhatted ones, the string field action for the non-BPS D-brane takes the same structural form as that in (2.8) and is given by
$$S=\frac{1}{4g^2}(e^{\widehat{\mathrm{\Phi }}}\widehat{Q}_Be^{\widehat{\mathrm{\Phi }}})(e^{\widehat{\mathrm{\Phi }}}\widehat{\eta }_0e^{\widehat{\mathrm{\Phi }}})_0^1𝑑t(e^{t\widehat{\mathrm{\Phi }}}_te^{t\widehat{\mathrm{\Phi }}})\{(e^{t\widehat{\mathrm{\Phi }}}\widehat{Q}_Be^{t\widehat{\mathrm{\Phi }}}),(e^{t\widehat{\mathrm{\Phi }}}\widehat{\eta }_0e^{t\widehat{\mathrm{\Phi }}})\},$$
(2.31)
where we have divided the overall normalization by a factor of two in order to compensate for the trace operation on the internal matrices. This action is invariant under the gauge transformation
$$\delta e^{\widehat{\mathrm{\Phi }}}=(\widehat{Q}_B\widehat{\mathrm{\Omega }})e^{\widehat{\mathrm{\Phi }}}+e^{\widehat{\mathrm{\Phi }}}(\widehat{\eta }_0\widehat{\mathrm{\Omega }}^{}),$$
(2.32)
where, as before the gauge transformation parameters $`\widehat{\mathrm{\Omega }}`$ and $`\widehat{\mathrm{\Omega }}^{}`$ are vertex operators with $`(n_g,n_p)`$ values $`(1,0)`$ and $`(1,1)`$ respectively. The internal CP indices carried by the gauge parameters are as follows
$$\widehat{\mathrm{\Omega }}=\mathrm{\Omega }_+\sigma _3+\mathrm{\Omega }_{}i\sigma _2,$$
(2.33)
with a similar relation holding for $`\widehat{\mathrm{\Omega }}^{}`$. The GSO even $`\mathrm{\Omega }_+`$ is Grassmann odd, while the GSO odd $`\mathrm{\Omega }_{}`$ is Grassmann even. This makes the overall gauge parameters $`\widehat{\mathrm{\Omega }}`$, $`\widehat{\mathrm{\Omega }}^{}`$ odd relative to $`\widehat{Q}_B`$, $`\widehat{\eta }_0`$. The proof of gauge invariance is formally identical to the one given in the earlier section. The equation of motion is just (2.23) with hats on fields and operators. Again, the gauge parameter $`\widehat{\mathrm{\Omega }}^{}`$ can be used to choose the gauge $`\xi _0\widehat{\mathrm{\Phi }}=0`$ so we can restrict to string states which are proportional to $`\xi _0`$.
For future use, we shall now give the expansion of the action (2.31) in power series in $`\widehat{\mathrm{\Phi }}`$. Expanding the exponentials in a power series, we get,
$`e^{\widehat{\mathrm{\Phi }}}𝒪e^{\widehat{\mathrm{\Phi }}}={\displaystyle \underset{M,N=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{(M+N+1)!}}\left({\displaystyle \genfrac{}{}{0pt}{}{M+N}{M}}\right)(1)^M\widehat{\mathrm{\Phi }}^M(𝒪\widehat{\mathrm{\Phi }})\widehat{\mathrm{\Phi }}^N,`$ (2.34)
valid for $`𝒪`$ equal to $`\widehat{Q}_B`$ or $`\widehat{\eta }_0`$. Using the cyclicity relations eq.(LABEL:e1aaa), and the identity (2.34), we can express the action (2.31) as
$$S=\frac{1}{2g^2}\underset{M,N=0}{\overset{\mathrm{}}{}}\frac{1}{(M+N+2)!}\left(\genfrac{}{}{0pt}{}{M+N}{N}\right)(1)^N(\widehat{Q}_B\widehat{\mathrm{\Phi }})\widehat{\mathrm{\Phi }}^M(\widehat{\eta }_0\widehat{\mathrm{\Phi }})\widehat{\mathrm{\Phi }}^N.$$
(2.35)
As in refs. we shall find it convenient to take the time direction to be periodic with period 1, so that for a static configuration we can identify the potential with the negative of the action. In this case, an analysis analogous to that in ref. shows that the string field action (2.31) describes a D-brane with mass
$$M=\frac{1}{2\pi ^2g^2}.$$
(2.36)
The details of this calculation have been outlined in appendix C. We shall calculate the tachyon potential and attempt to show that at the minimum it exactly cancels the mass $`M`$ given in eq.(2.36).
### 2.3 Superstring field theory on a D-brane anti-D-brane pair
This system incorporates a GSO$`(+)`$ sector in the form of vertex operators that represent strings that live on the brane or on the antibrane. With conventional Chan-Paton indices these would use the $`2\times 2`$ matrices $`I`$ and $`\sigma _3`$. In addition there is a GSO$`()`$ sector representing strings stretched between the brane and antibrane. With conventional Chan-Paton indices these would use the $`2\times 2`$ matrices $`\sigma _1`$ and $`\sigma _2`$. We will call the conventional Chan-Paton matrices external CP matrices, to distinguish them from the internal CP matrices used in the previous subsection. Since we still have the complication of including two GSO types in the string field, we will not dispense of the internal CP matrices, and thus the brane-antibrane system will use both internal and external CP matrices. This time we therefore write:
$$\widehat{\mathrm{\Phi }}=\mathrm{\Phi }_+^{(1)}II+\mathrm{\Phi }_+^{(2)}I\sigma _3+\mathrm{\Phi }_{}^{(3)}\sigma _1\sigma _1+\mathrm{\Phi }_{}^{(4)}\sigma _1\sigma _2,$$
(2.37)
where the first set of matrices are the internal ones and the second set are the external ones. In computing products of fields the two sets of matrices are defined to commute. For the operators and gauge parameters we have
$$\widehat{Q}_B=Q_B\sigma _3I,\widehat{\eta }_0=\eta _0\sigma _3I,$$
(2.38)
$$\widehat{\mathrm{\Omega }}=\mathrm{\Omega }_+^{(1)}\sigma _3I+\mathrm{\Omega }_+^{(2)}\sigma _3\sigma _3+\mathrm{\Omega }_{}^{(3)}i\sigma _2\sigma _1+\mathrm{\Omega }_{}^{(4)}i\sigma _2\sigma _2.$$
(2.39)
We are still writing all fields and operators with hats, for simplicity. The structure found earlier (eqs. (LABEL:e1aaa)-(2.30), in particular) survives when the correlators $`\mathrm{}`$ now include the double trace $`TrTr`$. The action takes then the same form as in (2.31) with the same normalization factor. If we restrict $`\widehat{\mathrm{\Phi }}`$ to be of the form $`\mathrm{\Phi }I\left(\begin{array}{cc}1& 0\\ 0& 0\end{array}\right)`$, we recover the open string field theory action (2.8) on a single BPS D-brane.
As discussed elsewhere, in analyzing the tachyon potential we can restrict ourselves to the external CP sector $`I`$ in the GSO$`(+)`$ sector, and to the external CP sector $`\sigma _1`$ in the GSO$`()`$ sector. Thus it is clear that there is a one to one correspondence between the component fields of the open string field theory on the non-BPS brane and that on the brane anti-brane system in this restricted sector. In fact since GSO$`()`$ fields must appear always in even numbers, the external CP factors with their trace will simply produce an extra factor of two for every interaction. Thus the computations of the tachyon potential are identical. For the same value of the open string coupling constant $`g^2`$, in the brane anti-brane system the potential is twice as large compared to that of the non-BPS D-brane due to the trace over the external CP factors. On the other hand now the mass of the brane or the anti-brane is given by $`(1/2\pi ^2g^2)`$, so that the total mass of the brane-antibrane system is also twice the mass of the non-BPS D-brane. Thus if for a non-BPS brane the potential energy at the bottom of the well cancels the tension, then for the brane-antibrane system the potential energy at the bottom of the well will also cancel the total tension of the brane-antibrane system. We will therefore use in this paper the simpler notation required for the analysis of the non-BPS brane.
## 3 Computation and analysis of the tachyon potential
We shall be interested here in the phenomenon of tachyon condensation on the non-BPS D-brane. As mentioned before, the analysis applies also to the D-brane anti-D-brane problem. We begin by setting up the level expansion of the full tachyon string field relevant to the condensation. We then discuss the expansion of the action. Finally, relegating some computations to an appendix, we calculate the tachyon potential, find its minimum and test the brane annihilation conjecture.
### 3.1 The tachyon string field
In the present case the zero momentum tachyon corresponds to the vertex operator $`\xi ce^\varphi \sigma _1`$. Let us denote by $`_1`$ the subset of vertex operators of ghost number 0 and picture number $`0`$, created from the matter superstress tensor $`(G_m(z),T_m(z))`$,<sup>3</sup><sup>3</sup>3For convenience of notation, we shall denote the $`k`$th oscillator mode of $`G_m`$ and $`T_m`$ by $`G_k^m`$ and $`L_k^m`$ respectively. and the ghost fields $`b`$, $`c`$, $`\xi `$, $`\eta `$, $`\varphi `$. It can be easily seen that by restricting the string field $`\widehat{\mathrm{\Phi }}`$ to be in $`_1`$ gives a consistent truncation of the action, and hence we can look for a solution of the equations of motion, representing tachyon condensation, by restricting the string field $`\widehat{\mathrm{\Phi }}`$ to be in this subspace $`_1`$. Thus from now on we shall always take the string field to lie in this restricted subspace.
We now expand the string field $`\widehat{\mathrm{\Phi }}`$ in a basis of $`L_0`$ eigenstates, and write the action (2.35) in terms of component fields, which are the coefficients of expansion of the string field in this basis. As in , we shall define the level of a string field component multiplying a vertex operator of conformal weight $`h`$ to be $`(h+\frac{1}{2})`$, so that the tachyon field, multiplying the vertex operator $`\xi ce^\varphi \sigma _1`$, has level 0. We also define the level of a given term in the string field action to be the sum of the levels of the individual fields appearing in that term, and define a level $`2n`$ approximation to the action to be the one obtained by including fields up to level $`n`$ and terms in the action up to level $`2n`$. Thus for example, a level 3 approximation to the action will involve string field components up to level (3/2). This is the approximation we shall be using to compute the action.
Using gauge invariance (2.32) of string field theory action, we can choose gauge conditions
$$b_0\widehat{\mathrm{\Phi }}=0,\xi _0\widehat{\mathrm{\Phi }}=0.$$
(3.2)
As in ref., the legitimacy of this gauge condition can be proved at the linearized level. We then assume that string field configuration under consideration is not too large, so that such a gauge choice is also possible for the configuration under study. Also, as discussed in ref., the gauge choice $`b_0\widehat{\mathrm{\Phi }}=0`$ can be made only for states with non-zero $`L_0`$ eigenvalue.
We can build systematically the relevant expansion of the string field by recalling that the string field $`\widehat{\mathrm{\Phi }}`$ satisfying the gauge condition $`\xi _0\widehat{\mathrm{\Phi }}=0`$ is related to the NS string field $`\widehat{V}`$ of by the relation $`\widehat{\mathrm{\Phi }}=\xi _0\widehat{V}`$. The string field $`\widehat{V}`$ is built on the tachyon vacuum $`|\stackrel{~}{\mathrm{\Omega }}c_1e^{\varphi (0)}|0`$. This vacuum state is GSO odd, it has ghost number $`+1`$ and $`L_0=1/2`$. Being in the minus one picture, it is annihilated by all positively moded oscillators $`\{\gamma _r,\beta _r\}`$. In addition, it is annihilated by $`b_0,L_1^m`$ and $`G_{\frac{1}{2}}^m`$. All relevant states of ghost number one are now obtained by acting with ghost number zero combinations of oscillators $`\{b,c,\beta ,\gamma ,L^m,G^m\}`$ on $`|\stackrel{~}{\mathrm{\Omega }}`$. The $`b_0\widehat{\mathrm{\Phi }}=0`$ gauge condition allows us to ignore states with a $`c_0`$ oscillator in them. The states one finds up to $`L_0`$ eigenvalue 1 are given in Table 1. For ease of notation we have not included the CP factor. Note that we have included at $`L_0=0`$ a state which is not annihilated by $`b_0`$. This is the case because having $`L_0=0`$ this state cannot be gauged away.
The string field we need, which uses the “large” Hilbert space, is obtained by acting on the states of the table with $`\xi _0`$. This operation, however, does not change the dimension of the operators. As shown in appendix B, the string field theory action in the restricted subspace $`_1`$ has a $`Z_2`$ twist symmetry under which string field components associated with a vertex operator of dimension $`h`$ carry charge $`(1)^{h+1}`$ for even $`2h`$, and $`(1)^{h+\frac{1}{2}}`$ for odd $`2h`$. The tachyon vertex operator, having dimension $`\frac{1}{2}`$, is even under this twist transformation. Thus we can consider a further truncation of the string field theory by restricting the string field $`\widehat{\mathrm{\Phi }}`$ to be twist even. This, in particular, means that the $`L_0=0`$ vertex operator mentioned above is to be omitted from the the string field. The same is true for the $`L_0=+1/2`$ state in the GSO$`()`$ sector. Therefore, in addition to the tachyon, we will include the three scalar fields appearing in the GSO$`(+)`$ sector at level $`3/2`$.
In the language of vertex operators $`|\stackrel{~}{\mathrm{\Omega }}`$ is $`ce^\varphi `$, and the three states in table 1 at level (3/2) are
$$c^2c\xi e^{2\varphi },\eta ,G_mce^\varphi ,$$
(3.3)
as can be seen with the help of eq.(2.1). We readily pass to the string field $`\widehat{\mathrm{\Phi }}`$ by acting the above operators with $`\xi _0`$, thus guaranteeing that both gauge conditions (3.2) are satisfied. Denoting the tachyon operator by $`\widehat{T}`$ and the three other operators by $`\widehat{A},\widehat{E}`$ and $`\widehat{F}`$ respectively, we have:
$`\widehat{T}`$ $`=\xi ce^\varphi \sigma _1`$ (3.4)
$`\widehat{A}`$ $`=c^2c\xi \xi e^{2\varphi }I`$ (3.5)
$`\widehat{E}`$ $`=\xi \eta I`$ (3.6)
$`\widehat{F}`$ $`=\xi G_mce^\varphi I`$ (3.7)
Therefore, the general twist even string field up to level $`(3/2)`$, satisfying the gauge condition (3.2), has the following form:<sup>4</sup><sup>4</sup>4Ref. had a factor of $`i`$ in front of the $`t\widehat{T}`$ term. In this paper we have used slightly different set of conformal maps $`f_k^{(N)}`$ in defining the string field theory action; these map the upper half plane into the inside of the unit disk rather than outside. With this choice, the kinetic term for the tachyon field has the standard sign provided there is no factor of $`i`$ multiplying $`t\widehat{T}`$ in eq.(3.8).
$$\widehat{\mathrm{\Phi }}=t\widehat{T}+a\widehat{A}+e\widehat{E}+f\widehat{F}.$$
(3.8)
As explained above, the tachyon vertex operator $`\widehat{T}`$ of $`L_0=1/2`$ is a GSO odd operator of level zero. The operators $`\widehat{A},\widehat{E}`$ and $`\widehat{F}`$ of $`L_0=+1`$, and thus level $`3/2`$, are in the GSO even sector.
### 3.2 Level expansion of the string action
We shall now substitute (3.8) into the action (2.35) and keep terms to all orders in $`t`$, but only up to quadratic order in $`a`$, $`e`$ and $`f`$. Although the string field action contains vertices of arbitrarily high order, it can be shown that the truncated action to any given level only has a finite number of terms. To see this, let us first note that for a term in the action of a given level, the number of fields of level $`>`$ 0 must be finite. Since all components of the string field other than the tachyon $`t`$ has level $`>`$ 0, we only need to show that there cannot be arbitrarily large number of tachyon fields. This is easily seen by noting that the tachyon vertex operator $`\widehat{T}`$ has $`1`$ unit of $`\varphi `$ momentum. Since in order to get a non-zero correlation function, the total $`\varphi `$ momentum of all the vertex operators must add up to $`2`$, it is clear that for a fixed set of other vertex operators, we can only insert a finite number of tachyon vertex operators in order to have a non-vanishing correlation function.
Each term in the action has one $`\widehat{\eta }_0`$ and one $`\widehat{Q}_B`$, each acting on a string field. While $`\widehat{\eta }_0`$ carries no $`\varphi `$-momentum, the BRST operator $`\widehat{Q}_B`$ can supply zero, one or two units of $`\varphi `$-momentum (see (2.4)). The operators $`\widehat{A}`$, $`\widehat{E}`$ and $`\widehat{F}`$ carry $`2`$, 0 and $`1`$ units of $`\varphi `$ momentum respectively. Since the operator $`E`$ entering in the string field carries no $`\varphi `$ momentum, this is the field that can appear together with the largest number of tachyon fields. For example, the string action term coupling $`E`$ with four $`T`$’s is nonvanishing since the tachyons give $`(4)`$ units of $`\varphi `$-momentum and the BRST operator can supply $`(+2)`$ units. Since we are going to compute the action to level three we can have a term in the string action with two $`E`$’s and four $`T`$’s. This is the term with the largest possible number of fields that can contribute to level three. This means we need the expansion of the string action (2.35) up to terms with six string fields. This is given by,
$`S`$ $`={\displaystyle \frac{1}{2g^2}}{\displaystyle \frac{1}{2}}(\widehat{Q}_B\widehat{\mathrm{\Phi }})(\widehat{\eta }_0\widehat{\mathrm{\Phi }})+{\displaystyle \frac{1}{6}}(\widehat{Q}_B\widehat{\mathrm{\Phi }})(\widehat{\mathrm{\Phi }}(\widehat{\eta }_0\widehat{\mathrm{\Phi }})(\widehat{\eta }_0\widehat{\mathrm{\Phi }})\widehat{\mathrm{\Phi }})`$
$`+{\displaystyle \frac{1}{24}}(\widehat{Q}_B\widehat{\mathrm{\Phi }})\left(\widehat{\mathrm{\Phi }}^2(\widehat{\eta }_0\widehat{\mathrm{\Phi }})2\widehat{\mathrm{\Phi }}(\widehat{\eta }_0\widehat{\mathrm{\Phi }})\widehat{\mathrm{\Phi }}+(\widehat{\eta }_0\widehat{\mathrm{\Phi }})\widehat{\mathrm{\Phi }}^2\right)`$
$`+{\displaystyle \frac{1}{120}}(\widehat{Q}_B\widehat{\mathrm{\Phi }})\left(\widehat{\mathrm{\Phi }}^3(\widehat{\eta }_0\widehat{\mathrm{\Phi }})3\widehat{\mathrm{\Phi }}^2(\widehat{\eta }_0\widehat{\mathrm{\Phi }})\widehat{\mathrm{\Phi }}+3\widehat{\mathrm{\Phi }}(\widehat{\eta }_0\widehat{\mathrm{\Phi }})\widehat{\mathrm{\Phi }}^2(\widehat{\eta }_0\widehat{\mathrm{\Phi }})\widehat{\mathrm{\Phi }}^3\right)`$
$`+{\displaystyle \frac{1}{720}}(\widehat{Q}_B\widehat{\mathrm{\Phi }})(\widehat{\mathrm{\Phi }}^4(\widehat{\eta }_0\widehat{\mathrm{\Phi }})4\widehat{\mathrm{\Phi }}^3(\widehat{\eta }_0\widehat{\mathrm{\Phi }})\widehat{\mathrm{\Phi }}+6\widehat{\mathrm{\Phi }}^2(\widehat{\eta }_0\widehat{\mathrm{\Phi }})\widehat{\mathrm{\Phi }}^24\widehat{\mathrm{\Phi }}(\widehat{\eta }_0\widehat{\mathrm{\Phi }})\widehat{\mathrm{\Phi }}^2+(\widehat{\eta }_0\widehat{\mathrm{\Phi }})\widehat{\mathrm{\Phi }}^4)`$
Since our computation will be restricted to twist even fields in $`_1`$, the above result can be simplified further by use of (B.16) and cyclicity. We find:
$`S`$ $`={\displaystyle \frac{1}{2g^2}}{\displaystyle \frac{1}{2}}(\widehat{Q}_B\widehat{\mathrm{\Phi }})(\widehat{\eta }_0\widehat{\mathrm{\Phi }})+{\displaystyle \frac{1}{3}}(\widehat{Q}_B\widehat{\mathrm{\Phi }})\widehat{\mathrm{\Phi }}(\widehat{\eta }_0\widehat{\mathrm{\Phi }})+{\displaystyle \frac{1}{12}}(\widehat{Q}_B\widehat{\mathrm{\Phi }})(\widehat{\mathrm{\Phi }}^2(\widehat{\eta }_0\widehat{\mathrm{\Phi }})\widehat{\mathrm{\Phi }}(\widehat{\eta }_0\widehat{\mathrm{\Phi }})\widehat{\mathrm{\Phi }})`$ (3.15)
$`+{\displaystyle \frac{1}{60}}(\widehat{Q}_B\widehat{\mathrm{\Phi }})\left(\widehat{\mathrm{\Phi }}^3(\widehat{\eta }_0\widehat{\mathrm{\Phi }})3\widehat{\mathrm{\Phi }}^2(\widehat{\eta }_0\widehat{\mathrm{\Phi }})\widehat{\mathrm{\Phi }}\right)`$
$`+{\displaystyle \frac{1}{360}}(\widehat{Q}_B\widehat{\mathrm{\Phi }})(\widehat{\mathrm{\Phi }}^4(\widehat{\eta }_0\widehat{\mathrm{\Phi }})4\widehat{\mathrm{\Phi }}^3(\widehat{\eta }_0\widehat{\mathrm{\Phi }})\widehat{\mathrm{\Phi }}+3\widehat{\mathrm{\Phi }}^2(\widehat{\eta }_0\widehat{\mathrm{\Phi }})\widehat{\mathrm{\Phi }}^2).`$
This expansion suffices for the present computation. All we need to do is to substitute (3.8) into this expression and evaluate the correlation functions appearing in various terms. The required calculations are relatively straightforward, and involve correlation functions of appropriate conformal transforms of the operators $`\widehat{T}`$, $`\widehat{A}`$, $`\widehat{E}`$, $`\widehat{F}`$, the BRST current $`j_B`$, and the field $`\eta `$. Here we shall only state the result; some of the details have been discussed in appendix D.<sup>5</sup><sup>5</sup>5We used the symbolic manipulation program Mathematica to carry out some of these computations.
### 3.3 The tachyon potential
We shall now give the result for the action and the potential by truncating it to level 3. Not all terms allowed by level counting are non-vanishing. Several vanish because they fail to satisfy $`\varphi `$-momentum conservation, for example, there is no $`a^2t^2`$ term.
The result, with $`S_k`$ denoting the level $`k`$ terms in the action, is
$`g^2S_0`$ $`={\displaystyle \frac{1}{4}}t^2{\displaystyle \frac{1}{2}}t^4,`$ (3.16)
$`g^2S_{\frac{3}{2}}`$ $`=at^2+{\displaystyle \frac{1}{4}}et^2+{\displaystyle \frac{5}{96}}\sqrt{50+22\sqrt{5}}et^4,`$ (3.18)
$`g^2S_3`$ $`=2ae5f^2`$ (3.24)
$`+\left({\displaystyle \frac{1}{\sqrt{2}}}{\displaystyle \frac{1}{24}}\right)e^2t^2{\displaystyle \frac{5}{18}}e^2t^4`$
$`{\displaystyle \frac{5}{4}}(4\sqrt{2}1)f^2t^2{\displaystyle \frac{1}{12}}(\mathrm{\hspace{0.17em}3}+40\sqrt{2})aet^2+{\displaystyle \frac{5}{12}}(\mathrm{\hspace{0.17em}10}\sqrt{2}1)eft^2.`$
The action to level three is given by $`S^{(3)}=S_0+S_{\frac{3}{2}}+S_3`$. Collecting the terms above, using the relation (2.36), and expressing the various radicals as approximate decimals, we have,
$`V=S`$ $`=`$ $`M(2\pi ^2)(0.25t^20.5t^4+at^2+0.25et^2+0.519et^4`$
$`2ae5f^2+0.665e^2t^20.278e^2t^4`$
$`5.82f^2t^24.96aet^2+\mathrm{\hspace{0.17em}\hspace{0.17em}5.476}eft^2)`$
The potential has extrema at $`(\pm t_0,a_0,e_0,f_0)`$ with
$$t_0=0.58882,a_0=0.056363,e_0=0.093175,f_0=0.012603.$$
(3.27)
At these extrema,
$$V=0.85446M.$$
(3.28)
The expected exact answer for the value of the potential at the extrema is $`M`$, so that it can cancel the mass of the D-brane exactly. Thus we see that the level three approximation produces 85% of the exact answer. Note that one finds 60% of the exact answer at level zero , so the approximation scheme appears to be converging to the exact answer.
The potential computed to this approximation (level three) includes the fields $`(a,e,f)`$ only quadratically. So they can be integrated out exactly to find an effective potential $`V(t)`$ for the tachyon. One obtains
$$v(t)\frac{1}{M}V(t)=4.93t^2\frac{(1+4.63t^2+3.21t^49.48t^611.67t^8)}{(1+1.16t^2)(1+2.48t^2)^2}$$
(3.29)
Several important properties are manifest from this expression. Since the denominators never vanish, there is no singularity in $`V(t)`$ to this approximation. The small tachyon instability for $`t0`$ is manifest. Since $`V+t^4`$, when $`t`$ is large, this potential is clearly bounded below. It can be easily checked that only critical points are $`\pm t_0`$ with $`t_0=0.58882`$; they are (equivalent) global minima of the presently computed effective potential.<sup>6</sup><sup>6</sup>6The critical point, however, is not a global minima of the full multiscalar potential $`V(t,a,e,f)`$. This is a reflection of the fact that some of the fields $`a,e,f`$ are auxiliary fields. The effective tachyon potential $`v(t)`$ has been displayed in fig. 1.
## 4 Tachyonic kink configuration
In the last section we analysed the tachyon potential on a non-BPS D-brane in a background independent fashion. In this section we focus on a specific case where the non-BPS D-brane corresponds to a non-BPS D-string of type IIA string theory wrapped on a circle of radius $`R`$. If $`𝒯_1`$ denotes the tension of the non-BPS D-string, then $`M=2\pi R𝒯_1`$. We denote by $`t`$ the (1+1) dimensional tachyon field living on the non-BPS D-string. Using the definition of $`v(t)`$ given in (3.29), we may express the potential $`V(t)`$ as $`2\pi R𝒯_1v(t)=𝒯_1_0^{2\pi R}𝑑xv(t)`$. From this we arrive at the conclusion that the potential energy of the D-string is given by
$$V(t)=𝒯_1_0^{2\pi R}𝑑xv(t).$$
(4.1)
We denote by $`x^\mu `$ for $`\mu =0,1`$ the world volume coordinates of the D-string, and $`xx^1`$ is the spatial coordinte.
Since $`v(t)`$ given in (3.29) has doubly degenerate minima at $`\pm t_0`$, we can consider a kink solution which interpolates between these two minima. It has been conjectured that this kink describes a BPS D0-brane of type IIA string theory. In this section we shall compute the mass of this kink solution, and compare it with the mass of a BPS D0-brane of type IIA string theory.
Computing the mass of the kink also requires knowledge of the kinetic term of the tachyon.<sup>7</sup><sup>7</sup>7By kinetic term we refer to the terms involving derivatives of the tachyon field $`t`$, including spatial derivatives. We use the kinetic term obtained from the quadratic term of the action, ignoring corrections from the higher order terms in the presence of background $`t`$. There is no justification for ignoring these corrections, and we should not expect this calculation to yield more than an order of magnitude estimate. Using the action (2.31), eq.(2.36) relating $`M=2\pi R𝒯_1`$ and $`g^2`$, and the fact that the length of D-string is $`2\pi R`$, we see that the Lagrangian contains a term
$$\frac{1}{2}(2\pi ^2𝒯_1)_0^{2\pi R}𝑑x_\mu t^\mu t.$$
(4.2)
We shall now take the limit $`R\mathrm{}`$, i.e. we consider infinitely long D-string. From eqs.(4.1) and (4.2) we get the following equations of motion for a static tachyonic configuration:
$$2\pi ^2_x^2t=v^{}(t).$$
(4.3)
Using the boundary condition that as $`x\pm \mathrm{}`$, $`t\pm t_0`$ and $`_xt0`$, the solution to this equation is implicitly given by:
$`_xt`$ $`=`$ $`{\displaystyle \frac{1}{\pi }}\sqrt{v(t)v(t_0)},`$
$`x`$ $`=`$ $`\pi {\displaystyle _0^{t(x)}}𝑑y{\displaystyle \frac{1}{\sqrt{v(y)v(t_0)}}}.`$ (4.4)
The total energy associated with this solution (measured above the $`t=t_0`$ solution), obtained by adding the kinetic and the potential terms is given by:
$$E=2\pi 𝒯_1_{t_0}^{t_0}𝑑y\sqrt{v(y)v(t_0)}.$$
(4.5)
Let $`𝒯_0`$ denote the mass of a BPS D0-brane of type IIA string theory. Then we have the relation:
$$𝒯_1=\sqrt{2}\frac{𝒯_0}{2\pi }.$$
(4.6)
Using this, eq.(4.5) can be written as
$$E=\sqrt{2}𝒯_0_{t_0}^{t_0}𝑑y\sqrt{v(y)v(t_0)}.$$
(4.7)
If we use the zeroeth order approximation for the potential
$$v(t)=2\pi ^2(\frac{1}{4}t^2+\frac{1}{2}t^4),$$
(4.8)
then (4.7) can be evaluated analytically, and we get
$$E=\frac{1}{6}\sqrt{2}\pi 𝒯_0.$$
(4.9)
This is about 74% of the expected answer $`𝒯_0`$.<sup>8</sup><sup>8</sup>8A similar result was independently obtained by Bergman, and also by Iqbal and Naqvi.
For the potential given in eq.(3.29), we can calculate the right hand side of eq.(4.7) numerically. The answer is
$$E=1.03𝒯_0.$$
(4.10)
Considering the crude approximation that we have used, this close agreement with the expected answer is likely to be accidental. However it is encouraging to note that the numerical answer is close to the expected answer.
This analysis can be easily extended to the case of a tachyonic kink on a non-BPS D-$`p`$ brane for any value of $`p`$.
## 5 Concluding remarks and open questions
There are two main points to the present paper. Point one: we seem to have a consistent NS open string field theory in which calculations are feasible. Point two: the calculations performed here with this string field theory give good direct evidence for the tachyon condensation phenomenon and its implications for unstable non-BPS D-branes as well as for the D-brane anti- D-brane system.
Let us first focus on the string field theory itself. While the cubic open string field theory of gives a consistent classical theory of bosonic open strings, its extension to superstrings was recognized early on to be problematic . Problems arise because the NS string vertex carries a picture changing operator at the interaction point, and in testing the associativity of the star product one induces the collision of two picture changing operators, upon which a divergence is encountered. It is believed that contact terms with infinite coefficients must be added to the action to restore gauge invariance. One may wonder if these complications are just irrelevant to the problem of computing the tachyon potential. We are not optimistic on this point. Indeed, in this theory, the potential of the tachyon field alone is purely quadratic. The absence of a cubic term (because of $`()^F`$ conservation), and the absence of a quartic term (as the theory is cubic) imply that the potential for the tachyon field alone has no critical points. It would therefore be necessary for the interactions of the tachyon with the other scalars to generate stabilizing terms of the right magnitude. However, the experience in this paper, as well as that in open bosonic string theory indicate that massive fields rather than stabilizing the tachyon, tend to lower the critical point that is generated by the tachyon field alone. Given the uncertainty in such arguments, it would be desirable to carry out the direct computation of the tachyon potential in this cubic theory. Since this theory is expressed in the “small Hilbert space” the table given in section 3.1 lists the relevant states. Just as it was the case in our present work, we expect that a twist analysis will show that the three fields at level $`3/2`$ are the ones that must be used for a lowest level nontrivial computation.
On the other hand the WZW-like NS string field theory used here is free of divergences and its gauge invariance is manifest. Given that it seems now to provide a consistent framework for dealing with the tachyon potential in the relevant brane systems, much of our work here has focused in the detailed setup of the action for the non-BPS brane, as well as for the brane anti-brane systems. We have also given very explicit consideration to the cyclicity and twist properties of the string action, and we have explained in detail how to work out branch cuts for dealing with the fractional dimension operators of the GSO odd sector. While this string field theory is non-polynomial, the level expansion is workable and the higher interactions are relatively simple to compute since they do not involve integration over the moduli space of Riemann surfaces; they are finite contact interactions. In contrast to bosonic string field theory, where gauge invariance was directly related to the covering of the moduli space of Riemann surfaces (see, for example, and ), something different and subtle is going on here as moduli spaces would be covered without the help of the higher interactions.<sup>9</sup><sup>9</sup>9Other approaches have been suggested to deal with the difficulties of . One possibility is to use string fields of non-canonical picture number . Another possibility is to make the superstring theory non-polynomial in the same way as must be done to incorporate closed strings off-shell . In such approach the region of moduli space where the collision of picture changing operators happens is within the interaction terms, which could be modified to prevent such collisions. Since the interactions in such theory would not be of contact type the level approximation would appear to be difficult to implement.
Turning now to the tachyon conjectures, the results obtained here are consistent with convergence to the expected values. While the condensation of the tachyon field alone gave about 60% of the desired value, the first nontrivial correction computed here (level 3) gave about 60% of the remaining energy. It should not be very hard to use the setup of this paper to carry the computation to level four, and perhaps to automate the computation further to deal with higher levels. Further evidence of convergence would be desirable. It would also be of interest to investigate further the properties of the tachyonic kink solution representing a lower dimensional brane. It should be noted that the convergence of the level approximation scheme to the answers seems slower than in the case of the bosonic open string, where the tachyon field alone gave about 70% of the desired energy, and inclusion of two additional scalars gave 95% of the expected answer.
Of course, at a deeper level the most intriguing questions remain those that were already apparent in the bosonic case : (i) Is there a closed form solution for the tachyon condensate? and, (ii) what is the physics of the vacuum around the tachyon condensate? Insight into any of these two questions would open up exciting new possibilities.
Acknowledgements: N.B. would like to thank Oren Bergman and Ion Vancea for useful discussions, Caltech, Harvard University and Massachusetts Institute of Technology for their hospitality, and CNPq grant 300256/94-9 for partial financial support. A.S. would like to thank Caltech for hospitality during part of this work. B.Z. would like to thank N. Moeller and W. Taylor for useful discussions.
## Appendix A Cyclicity property of string amplitudes
In this appendix we shall prove eqs.(LABEL:e1aa), (LABEL:e1aaa). Since the trace over the Chan Paton matrices satisfy the cyclicity property without any extra sign, we can work with the unhatted vertex operators, and prove (LABEL:e1aa) and (LABEL:e1aaa) simultaneously. First we shall prove this for string fields belonging to the restricted subspace $`_1`$, and then indicate its generalization for general string fields.
The cyclicity properties of the conformal field theory correlation functions are analyzed by using the property:
$`Tf_i^{(n)}A`$ $`=`$ $`f_{i+1}^{(n)}A\text{for}1i(n1)`$
$`Tf_n^{(n)}A`$ $`=`$ $`T^nf_1ARf_1A,`$ (A.1)
for any vertex operator A. Here $`T(w)=e^{2\pi i/n}w`$, and $`R=T^n`$ denotes rotation by $`2\pi `$. While the transformation $`R`$ acts trivially on the complex plane, it must be viewed in general as the composition $`T^n`$ of $`n`$ transformations by $`T`$. Thus $`R`$ affects the transformation of fields with non-integer dimension. Since $`T`$ maps unit disk to itself in a one to one fashion, it corresponds to an SL(2,R) transformation. Using $`SL(2,R)`$ invariance of the correlation functions on the disk, we can write
$$(f_1^{(n)}A_1)\mathrm{}(f_{n1}^{(n)}A_{n1})(f_n^{(n)}\mathrm{\Phi })=(f_2^{(n)}A_1)\mathrm{}(f_n^{(n)}A_{n1})(Rf_1^{(n)}\mathrm{\Phi })$$
(A.2)
In the subspace $`_1`$, the conformal weight of $`\mathrm{\Phi }`$ is integer if $`\mathrm{\Phi }`$ is Grassmann even (GSO$`(+)`$), and half integer if it is Grassman odd (GSO$`()`$). Thus the transformation by $`R`$ gives a factor of $`1`$ if $`\mathrm{\Phi }`$ is Grassmann even, and $`1`$ if $`\mathrm{\Phi }`$ is Grassmann odd. As can be seen from eq.(2.3), the product of all the operators inside the correlation function must be Grassmann even in order to get a non-vanishing correlator. Thus we pick up a factor of 1 ($`1`$) in moving the $`Rf_1^{(n)}\mathrm{\Phi }`$ factor on the right hand side of eq.(A.2) to the first place if $`\mathrm{\Phi }`$ is Grassmann even (odd). Thus the right hand side of eq.(A.2) may be written as
$$(f_1^{(n)}\mathrm{\Phi })(f_2^{(n)}A_1)\mathrm{}(f_n^{(n)}A_{n1})$$
(A.3)
irrespective of whether $`\mathrm{\Phi }`$ is Grassmann even or Grassmann odd.
If we replace $`\mathrm{\Phi }`$ by $`(Q_B\mathrm{\Phi })`$ or $`(\eta _0\mathrm{\Phi })`$, eq.(A.2) still holds, and transformation by $`R`$ still gives a factor of 1 ($`1)`$ if $`\mathrm{\Phi }`$ is Grassmann even (Grassmann odd). But now, since $`(Q_B\mathrm{\Phi })`$ and $`(\eta _0\mathrm{\Phi })`$ have statistics opposite to that of $`\mathrm{\Phi }`$, we pick up a factor of $`1`$ ($`+1)`$ in moving the $`Rf_1^{(n)}(Q_B\mathrm{\Phi })`$ or $`Rf_1^{(n)}(\eta _0\mathrm{\Phi })`$ factor to the first place if $`\mathrm{\Phi }`$ is Grassmann even (odd). This gives,
$`(f_1^{(n)}A_1)\mathrm{}(f_{n1}^{(n)}A_{n1})(f_n^{(n)}(Q_B\mathrm{\Phi }))`$ $`=`$ $`(f_1^{(n)}(Q_B\mathrm{\Phi }))(f_2^{(n)}A_1)\mathrm{}(f_n^{(n)}A_{n1})`$
$`(f_1^{(n)}A_1)\mathrm{}(f_{n1}^{(n)}A_{n1})(f_n^{(n)}(\eta _0\mathrm{\Phi }))`$ $`=`$ $`(f_1^{(n)}(\eta _0\mathrm{\Phi }))(f_2^{(n)}A_1)\mathrm{}(f_n^{(n)}A_{n1})`$
This proves eqs.(LABEL:e1aa) and (LABEL:e1aaa).
The cyclicity rules derived above also hold for a general string field $`\widehat{\mathrm{\Phi }}`$ not necessarily inside $`_1`$, and are in fact needed for the proof of gauge invariance of the action. The proof of these relations for a general D-brane system, however, requires using appropriate cyclicity axioms for the correlation functions of a general boundary conformal field theory. In the present context this axiom states that if $`\mathrm{\Phi }`$ denotes a vertex operator of conformal weight $`h`$, then in moving $`Rf_1^{(n)}\mathrm{\Phi }`$ from the extreme right to the extreme left in the right hand side of eq.(A.2), we pick up a factor of $`e^{2\pi ih}`$. On the other hand, from eq.(2.13) one can easily check that $`R=T^n`$ acting on $`\mathrm{\Phi }`$ gives a factor of $`e^{2\pi ih}`$. Thus these two factors cancel each other, and we recover the cyclicity rules given in the first of eqs.(LABEL:e1aa), (LABEL:e1aaa) for a general string field component $`\mathrm{\Phi }`$ or $`\widehat{\mathrm{\Phi }}`$. The other two equations of (LABEL:e1aa), (LABEL:e1aaa) can be proved along similar lines.
## Appendix B Twist invariance of the restricted action
In this appendix we shall show that the superstring field theory action (2.31), or equivalently (2.35), has a $`Z_2`$ twist invariance when we restrict the string field $`\widehat{\mathrm{\Phi }}`$ to lie in the subspace $`_1`$ defined in section 3. Using the form (2.35) of the action and the cyclicity relations given in eq.(LABEL:e1aaa), it is easy to verify that a vertex with even number of string fields ($`(M+N)`$ even terms of eq.(2.35)) is odd under $`Q_B\eta _0`$, whereas a vertex with odd number of string fields is even under $`Q_B\eta _0`$. Let us now consider a typical pair of terms in the string field theory action:
$$\widehat{\mathrm{\Phi }}^{i1}(\widehat{Q}_B\widehat{\mathrm{\Phi }})\widehat{\mathrm{\Phi }}^{ni1}(\widehat{\eta }_0\widehat{\mathrm{\Phi }})+(1)^{n+1}\widehat{\mathrm{\Phi }}^{i1}(\widehat{\eta }_0\widehat{\mathrm{\Phi }})\widehat{\mathrm{\Phi }}^{ni1}(\widehat{Q}_B\widehat{\mathrm{\Phi }})$$
(B.1)
Let now $`\widehat{\mathrm{\Phi }}_1,\mathrm{}\widehat{\mathrm{\Phi }}_n`$ denote $`n`$ arbitrary components of the string field $`\widehat{\mathrm{\Phi }}`$. In the expansion of the string action, the first term of (B.1) will give rise to a term of the form
$$(I)\widehat{\mathrm{\Phi }}_1\mathrm{}\widehat{\mathrm{\Phi }}_{i1}(\widehat{Q}_B\widehat{\mathrm{\Phi }}_i)\widehat{\mathrm{\Phi }}_{i+1}\mathrm{}\widehat{\mathrm{\Phi }}_{n1}(\widehat{\eta }_0\widehat{\mathrm{\Phi }}_n),$$
(B.2)
while the second term in (B.1) will give rise to a term of the form
$`(II)(1)^{n+1}\widehat{\mathrm{\Phi }}_{i1}\mathrm{}\widehat{\mathrm{\Phi }}_2\widehat{\mathrm{\Phi }}_1(\widehat{\eta }_0\widehat{\mathrm{\Phi }}_n)\widehat{\mathrm{\Phi }}_{n1}\mathrm{}\widehat{\mathrm{\Phi }}_{i+1}(\widehat{Q}_B\widehat{\mathrm{\Phi }}_i).`$ (B.3)
In fact, when expanding the string field in arbitrary components all terms in the action arising from (B.1) can be paired just as $`(I)`$ and $`(II)`$. Note that up to a cyclic transformation, the order of inputs in $`(I)`$ and $`(II)`$ are precisely reversed (twisted). We will relate $`(I)`$ to $`(II)`$ up to a sign, and this relation will enable us to derive a selection rule based on twist.
Let $`M(z)=z`$, $`\stackrel{~}{I}(z)=(1/z)`$, and $`R=T^n`$ denote rotation by $`2\pi `$. $`R`$ leaves vertex operators with integral conformal weight unchanged, and changes the sign of the vertex operators with half-integral conformal weight. As in the case of $`f_k^{(N)}`$’s, the definition of $`M`$ and $`\stackrel{~}{I}`$ are not complete unless we specify how to choose the sign when these transformations act on an half integral weight vertex operator. We adopt the following convention. Acting on a primary field $`\phi `$ of weight $`h`$,
$$M\phi (z)=e^{i\pi h}\phi (z),\stackrel{~}{I}\phi (z)=(iz)^{2h}\phi (\frac{1}{z}).$$
(B.4)
Note that since $`h`$ is either an integer or a half-integer, $`(iz)^{2h}`$ is well defined. We can now verify the relations:
$`f_i^{(n)}M\phi =\stackrel{~}{I}f_{ni+2}^{(n)}\phi \text{for}ni2`$ (B.5)
$`f_1^{(n)}M\phi =\stackrel{~}{I}Rf_1^{(n)}\phi ,`$ (B.6)
where the second equation is clearly the natural generalization of the first once we note that $`f_{n+1}^{(n)}`$ is identified with $`Rf_1^{(n)}`$. Since secondary vertex operators are obtained from products of derivatives of primary vertex operators, these relations also hold if we replace $`\phi `$ by a secondary vertex operator.
We now consider $`(I)`$ which explicitly reads
$`(I)`$ $`=`$ $`Tr(f_1^{(n)}\widehat{\mathrm{\Phi }}_1)\mathrm{}(f_{i1}^{(n)}\widehat{\mathrm{\Phi }}_{i1})(f_i^{(n)}(\widehat{Q}_B\widehat{\mathrm{\Phi }}_i))(f_{i+1}^{(n)}\widehat{\mathrm{\Phi }}_{i+1})`$ (B.7)
$`\mathrm{}(f_{n1}^{(n)}\widehat{\mathrm{\Phi }}_{n1})(f_n^{(n)}(\widehat{\eta }_0\widehat{\mathrm{\Phi }}_n)),`$
where for simplicity we have omitted the zeroes from the arguments of $`\widehat{\mathrm{\Phi }}_i`$. Since $`M`$ preserves the origin of the coordinate system, using (B.4), we can replace each vertex operator $`\widehat{\mathrm{\Phi }}_i(0)`$ in the above correlator by $`e^{i\pi h_i}M\widehat{\mathrm{\Phi }}_i(0)`$, where $`h_i`$ is the conformal dimension of $`\widehat{\mathrm{\Phi }}_i`$. We then use (LABEL:e7) to bring (B.7) into the form:
$`(I)=(1)^{{\scriptscriptstyle h_i}}Tr(\stackrel{~}{I}Rf_1^{(n)}\widehat{\mathrm{\Phi }}_1)(\stackrel{~}{I}f_n^{(n)}\widehat{\mathrm{\Phi }}_2)\mathrm{}(\stackrel{~}{I}f_{ni+3}^{(n)}\widehat{\mathrm{\Phi }}_{i1})`$
$`(\stackrel{~}{I}f_{ni+2}^{(n)}(\widehat{Q}_B\widehat{\mathrm{\Phi }}_i))(\stackrel{~}{I}f_{ni+1}^{(n)}\widehat{\mathrm{\Phi }}_{i+1})\mathrm{}(\stackrel{~}{I}f_3^{(n)}\widehat{\mathrm{\Phi }}_{n1})(\stackrel{~}{I}f_2^{(n)}(\widehat{\eta }_0\widehat{\mathrm{\Phi }}_n)).`$ (B.8)
We now use the following results:
* In the restricted sector in which we are working, the SL(2,C) transformation $`\stackrel{~}{I}`$ is a symmetry of the correlation functions. Thus we can remove all factors of $`\stackrel{~}{I}`$ from eq.(LABEL:e8).
* Acting on $`\widehat{\mathrm{\Phi }}_1`$, $`R`$ gives a factor of $`(1)^{2h_1}`$.
* If we reverse the ordering of $`f_n^{(n)}\widehat{\mathrm{\Phi }}_2\mathrm{}f_{ni+2}^{(n)}(\widehat{Q}_B\widehat{\mathrm{\Phi }}_i)\mathrm{}f_2^{(n)}(\widehat{\eta }_0\widehat{\mathrm{\Phi }}_n)`$ in eq.(LABEL:e8), then we pick up a factor of $`(1)(1)^{n_o^{}(n_o^{}1)/2}`$, where $`n_o^{}`$ is the number of odd string fields in the set $`\widehat{\mathrm{\Phi }}_2,\mathrm{}\widehat{\mathrm{\Phi }}_n`$. The first minus sign in this expression comes from passing $`\widehat{Q}_B`$ through $`\widehat{\eta }_0`$; the other factor comes from passing the odd components of the string field through each other. Note that due to the internal CP matrices there is no extra sign in passing $`\widehat{Q}_B`$ or $`\widehat{\eta }_0`$ through $`\widehat{\mathrm{\Phi }}_j`$, irrespective of whether $`\widehat{\mathrm{\Phi }}_j`$ is Grassmann even or odd.
Thus (LABEL:e8) can be written as
$`(I)=(1)^{{\scriptscriptstyle h_i}}(1)^{1+2h_1+n_o^{}(n_o^{}1)/2}Tr(f_1^{(n)}\widehat{\mathrm{\Phi }}_1)(f_2^{(n)}(\widehat{\eta }_0\widehat{\mathrm{\Phi }}_n))(f_3^{(n)}\widehat{\mathrm{\Phi }}_{n1})`$ (B.9)
$`\mathrm{}(f_{ni+1}^{(n)}\widehat{\mathrm{\Phi }}_{i+1})(f_{ni+2}^{(n)}(\widehat{Q}_B\widehat{\mathrm{\Phi }}_i))(f_{ni+3}^{(n)}\widehat{\mathrm{\Phi }}_{i1})\mathrm{}(f_n^{(n)}\widehat{\mathrm{\Phi }}_2)`$
$`=`$ $`(1)^{{\scriptscriptstyle h_i}}(1)^{1+2h_1+n_o^{}(n_o^{}1)/2}\widehat{\mathrm{\Phi }}_1(\widehat{\eta }_0\widehat{\mathrm{\Phi }}_n)\widehat{\mathrm{\Phi }}_{n1}\mathrm{}\widehat{\mathrm{\Phi }}_{i+1}(\widehat{Q}_B\widehat{\mathrm{\Phi }}_i)\widehat{\mathrm{\Phi }}_{i1}\mathrm{}\widehat{\mathrm{\Phi }}_2.`$
Let $`n_e`$ and $`n_o`$ denote the total number of even and odd fields in the set $`\widehat{\mathrm{\Phi }}_1,\mathrm{}\widehat{\mathrm{\Phi }}_n`$. Since $`n_o`$ is always even, we may write $`n_o=2m`$ for some integer $`m`$. We now analyse two cases separately.
1. $`\widehat{\mathrm{\Phi }}_1`$ odd. In this case $`n_o^{}=2m1`$ and $`(1)^{2h_1}=1`$. Thus:
$$(1)^{2h_1}(1)^{n_o^{}(n_o^{}1)/2}=(1)^{1+(m1)(2m1)}=(1)^m,$$
(B.10)
for integer $`m`$.
2. $`\widehat{\mathrm{\Phi }}_1`$ is even. In this case $`n_o^{}=2m`$, $`(1)^{2h_1}=1`$, and we have
$$(1)^{2h_1}(1)^{n_o^{}(n_o^{}1)/2}=(1)^{m(2m1)}=(1)^m.$$
(B.11)
Thus in both cases $`(1)^{2h_1}(1)^{n_o^{}(n_o^{}1)/2}=(1)^m`$. Using eqs.(B.10)-(B.11), and the cyclicity property (LABEL:e1aaa), we can finally express the right hand side of (LABEL:e9) as
$$(I)=(1)^{\frac{n_o}{2}+1}(1)^{{\scriptscriptstyle h_i}}\widehat{\mathrm{\Phi }}_{i1}\mathrm{}\widehat{\mathrm{\Phi }}_2\widehat{\mathrm{\Phi }}_1(\widehat{\eta }_0\widehat{\mathrm{\Phi }}_n)\widehat{\mathrm{\Phi }}_{n1}\mathrm{}\widehat{\mathrm{\Phi }}_{i+1}(\widehat{Q}_B\widehat{\mathrm{\Phi }}_i).$$
(B.12)
We now recognize that the operators inside the correlator are ordered just as in (B.3). Since the total contribution to the action is given by the addition of $`(I)`$ and $`(II)`$, combining (B.12) and (B.3) we see that we get a non-zero contribution only if
$$(1)^{n+\frac{n_o}{2}}(1)^{_ih_i}=1.$$
(B.13)
Since $`n_o`$ is always even, we have $`(1)^n=(1)^{n_o+n_e}=(1)^{n_e}`$. Thus we may rewrite (B.13) as
$$(1)^{_{even}(h_i+1)}(1)^{_{odd}(h_i+\frac{1}{2})}=1.$$
(B.14)
This can be interpreted by saying that the action has a $`Z_2`$ twist invariance under which even fields carry twist charge $`(1)^{h+1}`$ and odd fields carry twist charge $`(1)^{h+\frac{1}{2}}`$. This means, in particular, that the tachyon $`\widehat{T}`$, being Grassmann odd and of dimension $`1/2`$ has twist charge $`+1`$. In the computation of the tachyon potential we can therefore restrict $`_1`$ to twist even fields.
The results of this appendix can be used to relate terms in the string action. It follows from our discussion above that for a vertex involving $`n`$ string fields in $`_1`$:
$$\widehat{\mathrm{\Phi }}_1\mathrm{}(\widehat{Q}_B\widehat{\mathrm{\Phi }}_k)\mathrm{}(\widehat{\eta }_0\widehat{\mathrm{\Phi }}_l)\mathrm{}\widehat{\mathrm{\Phi }}_n=()^{n+1}\left(\underset{i=1}{\overset{n}{}}\mathrm{\Omega }_i\right)\widehat{\mathrm{\Phi }}_n\mathrm{}(\widehat{\eta }_0\widehat{\mathrm{\Phi }}_l)\mathrm{}(\widehat{Q}_B\widehat{\mathrm{\Phi }}_k)\mathrm{}\widehat{\mathrm{\Phi }}_1$$
(B.15)
where $`\mathrm{\Omega }_i`$ is the twist eigenvalue of $`\widehat{\mathrm{\Phi }}_i`$. When we restrict to twist even fields in $`_1`$ the above equation is even simpler:
$$\widehat{\mathrm{\Phi }}_1\mathrm{}(\widehat{Q}_B\widehat{\mathrm{\Phi }}_k)\mathrm{}(\widehat{\eta }_0\widehat{\mathrm{\Phi }}_l)\mathrm{}\widehat{\mathrm{\Phi }}_n=()^{n+1}\widehat{\mathrm{\Phi }}_n\mathrm{}(\widehat{\eta }_0\widehat{\mathrm{\Phi }}_l)\mathrm{}(\widehat{Q}_B\widehat{\mathrm{\Phi }}_k)\mathrm{}\widehat{\mathrm{\Phi }}_1.$$
(B.16)
## Appendix C Mass of the D-brane
In this appendix we shall show that the mass of the D-brane, whose world volume theory is given by the action (2.31), is given by $`(2\pi ^2g^2)^1`$. The strategy that we shall be following is as follows. As in ref., we assume that there is a set of (at least one) non-compact flat directions transverse to the D-brane; we shall denote these coordinates by $`x^i`$. Then the open string modes living on the D-brane will include the location of the D-brane along the directions $`x^i`$. Let $`Y^i`$ denote the coordinate of the D-brane along $`x^i`$. The string field theory action contains a term proportional to $`(_tY^i)^2`$, where $`_t`$ denotes time derivative. The coefficient of the $`(_tY^i)^2`$ can be identified as half of the D-brane mass.
Let $`X^i`$ be the free world-volume scalar field associated with the coordinate $`x^i`$, and $`\psi ^i`$, $`\stackrel{~}{\psi }^i`$ its left- and right-handed supersymmetric partners. We denote by $`x^0t`$ the time coordinate, $`X^0`$ the corresponding world-volume scalar field, and $`k_0`$ the quantum number labelling momentum conjugate to $`X^0`$. If we write $`X^\mu =X_L^\mu +X_R^\mu `$ with $`L`$ and $`R`$ denoting left and right-moving components, then,
$$X_L^\mu (z)X_L^\nu (w)\frac{\eta ^{\mu \nu }}{2(zw)^2},\psi ^\mu (z)\psi ^\nu (w)\frac{\eta ^{\mu \nu }}{2(zw)},$$
(C.1)
with $`\eta ^{\mu \nu }=diag(1,1,\mathrm{},1)`$. With this normalization,
$$T_m=X_LX_L\psi \psi +\mathrm{},G_m=2i\psi X_L+\mathrm{}.$$
(C.2)
There is a similar set of relations for the right-moving (anti-holomorphic) fields.
Since the time direction has been taken to be periodic with period 1, $`k_0`$ is quantized in units of $`2\pi `$. Let us now consider the following term in the expansion of the string field $`\widehat{\mathrm{\Phi }}`$;
$$\widehat{\mathrm{\Phi }}=\underset{k_0}{}\varphi ^i(k_0)\sqrt{2}\xi ce^\varphi \psi ^ie^{ik_0X^0}I+\mathrm{}.$$
(C.3)
The $`\sqrt{2}`$ factor in this expansion has been included to compensate for the factor of (1/2) in the operator product of $`\psi ^i`$ with itself. Although $`X^0=X_L^0+X_R^0`$, using the Neumann boundary condition $`X_L=X_R`$ at the boundary we can replace $`e^{ik_0X^0}`$ by $`e^{2ik_0X_L^0}`$. This facilitates computation of various correlation functions. In particular, using eqs.(C.1), (C.2) we see that this vertex operator has $`L_0^m`$ eigenvalue equal to $`\frac{1}{2}(k_0)^2`$, where $`L_k^m`$ denotes the $`k`$th mode of the matter Virasoro generator.
We shall now examine the quadratic term in the action involving the mode $`\varphi ^i(k_0)`$. Only the $`cT_m`$ term of the BRST current $`j_B`$ contributes to the $`k_0`$ dependent part of the quadratic term involving this mode, and the result is given by
$$\frac{1}{2g^2}\underset{k_0}{}(k_0)^2\varphi ^i(k_0)\varphi ^i(k_0),$$
(C.4)
in the $`\alpha ^{}=1`$ unit. If $`\chi ^i(t)_{k_0}e^{ik_0t}\varphi ^i(k_0)`$ denotes the Fourier transform of $`\varphi ^i(k_0)`$, then the above action can be rewritten as
$$\frac{1}{2g^2}𝑑t_t\chi ^i_t\chi ^i,$$
(C.5)
where $`tx^0`$ denotes the time variable conjugate to $`k_0`$.
Up to an overall normalization factor, $`\chi ^i`$ has the interpretation of the location $`Y^i`$ of the D-brane in the $`x^i`$ direction. We shall now determine the normalization factor between $`\chi ^i`$ and $`Y^i`$. For this, instead of taking a single D-brane, let us take a pair of identical D-branes, separated by a distance $`b^i`$ along the $`X^i`$ direction. Then each state in the open string Hilbert space carries a $`2\times 2`$ Chan Paton factor, besides the usual CP factor carried by a single non-BPS D-brane; we shall call these external CP factors. States with off diagonal external CP factors, representing open strings stretched between the two branes, are forced to carry an amount of winding charge $`b^i`$ along $`X^i`$. For $`\alpha ^{}=1`$, i.e. string tension$`=(2\pi )^1`$, the classical contribution to the mass of these open string states due to the tension of the string is equal to $`|\stackrel{}{b}|/(2\pi )`$. If we now move one of the branes by an amount $`Y^i`$ along $`X^i`$, the change in the (mass)<sup>2</sup> of the open string with Chan Paton factors $`\left(\begin{array}{cc}0& 1\\ 0& 0\end{array}\right)`$ and $`\left(\begin{array}{cc}0& 0\\ 1& 0\end{array}\right)`$ should be given by:
$$\frac{1}{(2\pi )^2}\{(\stackrel{}{b}+\stackrel{}{Y})^2\stackrel{}{b}^2\}=\frac{1}{2\pi ^2}\stackrel{}{b}\stackrel{}{Y}+O(\stackrel{}{Y}^2).$$
(C.6)
On the other hand, since $`\chi ^i`$ denotes the mode which translates the brane, moving one of the branes along $`X^i`$ will correspond to switching on a constant $`\chi ^i`$. This is represented by a string field background
$$\sqrt{2}\chi ^i\xi ce^\varphi \psi ^iI\left(\begin{array}{cc}1& 0\\ 0& 0\end{array}\right).$$
(C.7)
We shall now explicitly use the string field theory action (2.31) to calculate the change of the (mass)<sup>2</sup> of states with Chan Paton factors $`\left(\begin{array}{cc}0& 1\\ 0& 0\end{array}\right)`$ and $`\left(\begin{array}{cc}0& 0\\ 1& 0\end{array}\right)`$ due to the presence of this background string field, and compare with eq.(C.6). For this we note that the vertex operator for the lowest mass open string with internal CP factor $`I`$ and external CP factors $`\left(\begin{array}{cc}0& 1\\ 0& 0\end{array}\right)`$ and $`\left(\begin{array}{cc}0& 0\\ 1& 0\end{array}\right)`$ are given by, respectively,
$`\xi ce^\varphi (\stackrel{}{ϵ}\stackrel{}{\psi })e^{i\frac{b^i}{2\pi }(X_L^iX_R^i)}e^{2ik_0X_L^0}I\left(\begin{array}{cc}0& 1\\ 0& 0\end{array}\right),\text{and}`$
$`\xi ce^\varphi (\stackrel{}{ϵ}\stackrel{}{\psi })e^{i\frac{b^i}{2\pi }(X_L^iX_R^i)}e^{2ik_0X_L^0}I\left(\begin{array}{cc}0& 0\\ 1& 0\end{array}\right),`$ (C.8)
where $`\stackrel{}{ϵ}`$ is a polarization vector. Using Dirichlet boundary condition on $`X^i`$, we can write $`X_L^iX_R^i=2X_L^i`$. Requiring BRST invariance of these vertex operators gives,
$$\stackrel{}{ϵ}\stackrel{}{b}=0,(k_0)^2=\frac{\stackrel{}{b}^2}{(2\pi )^2}.$$
(C.9)
Thus they represent states of mass $`|\stackrel{}{b}|/(2\pi )`$. We shall normalize $`\stackrel{}{ϵ}`$ such that
$$|\stackrel{}{ϵ}|^2=2.$$
(C.10)
Let us now consider the following expansion of the string field
$$\widehat{\mathrm{\Phi }}=\chi ^i\widehat{P}^i+\underset{k_0}{}(u(k_0)\widehat{U}(k_0)+u^{}(k_0)\widehat{V}(k_0))+\mathrm{}$$
(C.11)
where
$$\widehat{P}^i=\sqrt{2}\xi ce^\varphi \psi ^iI\left(\begin{array}{cc}1& 0\\ 0& 0\end{array}\right),$$
(C.12)
$$\widehat{U}(k_0)=\xi ce^\varphi (\stackrel{}{ϵ}\stackrel{}{\psi })e^{2i\frac{b^i}{2\pi }X_L^i}e^{2ik_0X_L^0}I\left(\begin{array}{cc}0& 1\\ 0& 0\end{array}\right),$$
(C.13)
$$\widehat{V}(k_0)=\xi ce^\varphi (\stackrel{}{ϵ}\stackrel{}{\psi })e^{2i\frac{b^i}{2\pi }X_L^i}e^{2ik_0X_L^0}I\left(\begin{array}{cc}0& 0\\ 1& 0\end{array}\right).$$
(C.14)
$`\chi ^i`$, $`u(k_0)`$ and $`u^{}(k_0)`$ are specific components of the string field. We can now evaluate the string field theory action as a function of these fields. We shall be interested in the quadratic term involving $`u`$, $`u^{}`$, as well as the $`\chi ^iuu^{}`$ coupling. The quadratic term is given by
$$\frac{1}{g^2}\underset{k_0}{}u^{}(k_0)u(k_0)(k_0^2\frac{\stackrel{}{b}^2}{4\pi ^2}).$$
(C.15)
The computation of the $`\chi ^iuu^{}`$ coupling can be simplified if we work on-shell at $`k_0^2=\stackrel{}{b}^2/(2\pi )^2`$. (This suffices for computing the shift in mass<sup>2</sup> of the state to order $`\chi ^i`$.) We now note that:
* Using the three point vertex $`(12g^2)^1((\widehat{Q}_B\widehat{\mathrm{\Phi }})\widehat{\mathrm{\Phi }}(\widehat{\eta }_0\widehat{\mathrm{\Phi }})\widehat{\mathrm{\Phi }}(\widehat{Q}_B\widehat{\mathrm{\Phi }})(\widehat{\eta }_0\widehat{\mathrm{\Phi }}))`$ we get twelve terms contributing to the $`\chi ^iuu^{}`$ coupling. Half of these terms vanish due to the trace identity:
$$Tr\left(\left(\begin{array}{cc}0& 0\\ 1& 0\end{array}\right)\left(\begin{array}{cc}0& 1\\ 0& 0\end{array}\right)\left(\begin{array}{cc}1& 0\\ 0& 0\end{array}\right)\right)=0.$$
(C.16)
The other cyclic ordering of these matrices produce a non-zero answer (equal to unity) for this trace.
* Each of the vertex operators $`\widehat{P}^i`$, $`\widehat{U}`$ and $`\widehat{V}`$ is annihilated by $`\widehat{Q}_B\widehat{\eta }_0`$ if $`k_0^2=\stackrel{}{b}^2/(2\pi )^2`$. Using this result we can manipulate each of the remaining six terms so that $`\widehat{Q}_B`$ acts on $`\widehat{P}^i`$, and $`\widehat{\eta }_0`$ acts on $`\widehat{U}`$. Finally, using the cyclicity relations (LABEL:e1aaa) we can show that each of these six terms gives identical result proportional to $`(\widehat{Q}_B\widehat{P}^i)(\widehat{\eta }_0\widehat{U})\widehat{V}`$.
After performing the trace over CP factors, and restricting to only on-shell components of $`u`$ and $`u^{}`$, we may express the $`\chi ^iuu^{}`$ term in the action as:
$$\frac{1}{g^2}\underset{k_0}{}^{}\chi ^iu^{}(k_0)u(k_0)f_1(Q_BP^i)f_2(\eta _0U(k_0))f_3V(k_0),$$
(C.17)
where $`^{}`$ denotes sum over on-shell momenta $`k_0=\pm |\stackrel{}{b}|/(2\pi )`$. This correlation function is easily evaluated and the result is
$$\frac{1}{g^2}\frac{1}{\sqrt{2}\pi }\stackrel{}{b}\stackrel{}{\chi }\underset{k_0}{}^{}u^{}(k_0)u(k_0).$$
(C.18)
Combining this with eq.(C.15) we see that the shift in the mass<sup>2</sup> of the $`u`$, $`u^{}`$ field due to the presence of $`\chi ^i`$ background is given by
$$\frac{1}{\sqrt{2}\pi }\stackrel{}{b}\stackrel{}{\chi }+O(\stackrel{}{\chi }^2).$$
(C.19)
Comparing eqs.(C.6) and (C.19) we get
$$\chi ^i=\frac{Y^i}{\sqrt{2}\pi }.$$
(C.20)
Once we have determined the relative normalization between $`\chi ^i`$ and $`Y^i`$, we can return to the system containing a single brane.<sup>10</sup><sup>10</sup>10This can be done, for example, by moving the other brane infinite distance away by taking the limit $`|\stackrel{}{b}|\mathrm{}`$. Substituting eq.(C.20) into eq.(C.5), we get,
$$(4\pi ^2g^2)^1𝑑t_tY^i_tY^i.$$
(C.21)
This contribution to the D-brane world-volume action can be interpreted as due to the kinetic energy associated with the collective motion of the D-brane in the non-compact transverse directions. This allows us to identify the D-brane mass as
$$M=(2\pi ^2g^2)^1.$$
(C.22)
## Appendix D Details on the calculation of the tachyon potential
We first consider some of the ingredients of the calculation, then do a particular example in detail. First of all, computation of $``$ involving the various vertex operators $`T,A,\mathrm{}`$ requires knowledge of $`fT(0)`$, $`fA(0)`$ etc., for a conformal map $`f`$. If $`f(0)=w`$, then we have the following relations:
$`fT(0)`$ $`=`$ $`(f^{}(0))^{\frac{1}{2}}T(w)`$ (D.1)
$`fA(0)`$ $`=`$ $`f^{}(0)\left(A(w){\displaystyle \frac{f^{\prime \prime }(0)}{(f^{}(0))^2}}cc\xi \xi e^{2\varphi }(w)\right)`$ (D.2)
$`fE(0)`$ $`=`$ $`f^{}(0)\left(E(w){\displaystyle \frac{f^{\prime \prime }(0)}{2(f^{}(0))^2}}\right)`$ (D.3)
$`fF(0)`$ $`=`$ $`f^{}(0)F(w)`$ (D.4)
Since the action involves $`Q_B`$ and $`\eta _0`$ acting on string fields, we need to evaluate those on $`T,A,E`$ and $`F`$ and the result of conformal transform of these operators. However the analysis can be simplified by noting that
$$f(𝒪A)=𝒪(fA),$$
(D.5)
where $`𝒪`$ can be either $`Q_B`$ or $`\eta _0`$. This is due to the fact that the BRST current $`j_B`$ and $`\eta `$ are dimension 1 primary fields. Thus for example, in calculating correlation function involving $`f(Q_BA(0))`$ we need to calculate the correlation function involving $`j_B(w)fA(0)`$ and pick up the residue of the pole at $`w=f(0)`$. A similar procedure holds for $`f(\eta _0A(0))`$.
These relations, together with eq.(2.6), and the identity
$`{\displaystyle \underset{i=1}{\overset{n+1}{}}}\xi (x_i){\displaystyle \underset{j=1}{\overset{n}{}}}\eta (y_j){\displaystyle \underset{k=1}{\overset{m}{}}}b(u_k){\displaystyle \underset{l=1}{\overset{m+3}{}}}c(v_l){\displaystyle \underset{s=1}{\overset{p}{}}}e^{q_s\varphi (z_s)}`$
$`=`$ $`{\displaystyle \underset{i<i^{}}{}}(x_ix_i^{}){\displaystyle \underset{j<j^{}}{}}(y_jy_j^{}){\displaystyle \underset{i,j}{}}(x_iy_j)^1{\displaystyle \underset{k<k^{}}{}}(u_ku_k^{}){\displaystyle \underset{l<l^{}}{}}(v_lv_l^{}){\displaystyle \underset{k,l}{}}(u_kv_l)^1`$ (D.6)
$`\times {\displaystyle \underset{s<s^{}}{}}(z_sz_s^{})^{q_sq_s^{}},`$
allows us to compute the relevant terms which appear in the computation of the tachyon potential. Eq.(LABEL:ecorrln) follows from the normalization convention (2.3), and the operator products (2.7).
In evaluating correlation functions involving the operator $`E`$, we need to exercise special care, as it involves product of $`\xi `$ and $`\eta `$ at the same point. This has to be interpreted as:
$$\xi \eta (w)=\underset{zw}{lim}\left(\xi (z)\eta (w)\frac{1}{zw}\right).$$
(D.7)
Let us give as an example the computation of the quartic term in the tachyon potential. From the expansion of the action (3.15), focusing on the terms with four string fields, we find:
$`g^2S|_{t^4}`$ $`={\displaystyle \frac{t^4}{24}}\left\{(\widehat{Q}_B\widehat{T})\widehat{T}(\widehat{\eta }_0\widehat{T})\widehat{T}(\widehat{Q}_B\widehat{T})\widehat{T}\widehat{T}(\widehat{\eta }_0\widehat{T})\right\},`$ (D.9)
$`={\displaystyle \frac{t^4}{12}}\left\{(Q_BT)T(\eta _0T)T+(Q_BT)TT(\eta _0T)\right\}.`$
In the second step we evaluated the trace over the internal CP matrices. We therefore have two correlators to compute. Using the fact that $`T`$ correspond to a dimension $`(1/2)`$ primary field, and that both $`j_B(w)`$ and $`\eta (w)`$ have only single poles near a $`T`$, the first correlator in the above equation can be written as:
$`C(f_1,f_2,f_3,f_4)`$ $``$ $`f_1(Q_BT(0))f_2T(0)f_3(\eta _0T(0))f_4T(0)`$
$`=`$ $`\underset{y_1w_1}{lim}\underset{y_2w_3}{lim}(y_1w_1)(y_2w_3){\displaystyle \frac{j_B(y_1)T(w_1)T(w_2)\eta (y_2)T(w_3)T(w_4)}{(f_1^{})^{\frac{1}{2}}(f_2^{})^{\frac{1}{2}}(f_3^{})^{\frac{1}{2}}(f_4^{})^{\frac{1}{2}}}},`$
where $`w_i=f_i(0)`$. We have, for simplicity of notation, defined $`f_if_i^{(4)}`$. This correlation function can be easily evaluated, and the answer is
$$C(f_1,f_2,f_3,f_4)=\frac{w_{13}w_{24}}{(f_1^{})^{\frac{1}{2}}(f_2^{})^{\frac{1}{2}}(f_3^{})^{\frac{1}{2}}(f_4^{})^{\frac{1}{2}}},$$
(D.13)
where $`w_{ij}=(w_iw_j)`$. We now recognize that the second correlator in (D.9) is simply $`C(f_1,f_2,f_4,f_3)`$ with no extra sign factor because the last two vertex operators do not induce a sign factor when they are transposed. We can therefore write the complete answer as
$$g^2S|_{t^4}=\frac{w_{13}w_{24}+w_{14}w_{23}}{12(f_1^{})^{\frac{1}{2}}(f_2^{})^{\frac{1}{2}}(f_3^{})^{\frac{1}{2}}(f_4^{})^{\frac{1}{2}}}t^4.$$
(D.14)
This off-shell amplitude is PSL(2,C) invariant<sup>11</sup><sup>11</sup>11See for a Riemann surface interpretation of invariant off-shell amplitudes.. Indeed letting
$$w\frac{aw+b}{cw+d},adbc=1,$$
(D.15)
we readily find that
$$w_{ij}\frac{w_{ij}}{(cw_i+d)(cw_j+d)},f_i^{}\frac{f_i^{}}{(cw_i+d)^2}$$
(D.16)
and therefore we get PSL(2,C) invariance if we choose the branch<sup>12</sup><sup>12</sup>12In the chosen SL(2,C) transformation there is a sign ambiguity in which all coefficients $`a,b,c,d`$ of the transformation are changed in sign. Since this transformation must be used for an even number of punctures, this is not a problem.
$$(f_i^{})^{1/2}\frac{(f_i^{})^{1/2}}{cw_i+d}.$$
(D.17)
We evaluate now the term. Our first choice of coordinates is that of the unit disk, described in detail in section 2.1. The prescription for dealing with the square roots there ((2.13)) is used to find
$$g^2S|_{t^4}=\frac{2(2i)+(1+i)^2}{12e^{i\pi /2}}t^4=\frac{1}{2}t^4$$
(D.18)
which is the result obtained in .
We shall now do the computation in the upper half plane (UHP) using the maps $`g_k^{(n)}`$, related to $`f_k^{(n)}`$ by an SL(2,C) transformation which maps the disk to UHP. But before we proceed, we need to derive the analog of eq.(2.13) for half integer $`h`$, i.e. the presription for choosing the sign of $`(g_k^{(n)}(0))^{\frac{1}{2}}`$ appearing in the conformal transform of half-integer weight fields. This will be done by starting with the presciption (2.13) and then using prescription (D.17) for an appropriate SL(2,C) transformation relating $`f_k^{(n)}`$ to $`g_k^{(n)}`$. First note that for fixed $`n`$ and $`k`$, $`f_k^{(n)}(z)`$ moves anti-clockwise along the boundary of the unit disk as $`z`$ moves along the positively oriented real line.<sup>13</sup><sup>13</sup>13This is related to the fact that the canonical half disks representing the strings are all mapped analytically into the interior of the disk and the boundary of the canonical half-disks is oriented in the direction of increasing real values. In addition, since the map from the disk to the UHP takes the anti-clockwise oriented boundary of the disk to the positively oriented real line, it is clear that the $`g_k^{(n)}`$’s map to positive real values at the punctures on the real line. In computing $`(g_k^{(n)})^{\frac{1}{2}}`$ we have a sign ambiguity. We shall now show that if the conformal map relating $`f_k^{(n)}`$ to $`g_k^{(n)}`$ is such that the points $`g_1^{(n)}(0),\mathrm{}g_n^{(n)}(0)`$ are ordered from the left to the right on the real axis, then we should choose the positive sign for all the $`(g_k^{(n)}(0))^{\frac{1}{2}}`$.
We prove this as follows. As a first step it is convenient to rotate the punctures on the disk to a new position. For this, we define:
$$\stackrel{~}{f}_k^{(n)}(z)=e^{\frac{2\pi i}{n}(\frac{n}{2}+1ϵ)}f_k^{(n)}(z)=e^{\frac{2\pi i}{n}(k\frac{n}{2}ϵ)}\left(\frac{1+iz}{1iz}\right)^{\frac{2}{n}},$$
(D.19)
where $`ϵ`$ is a small positive number; in fact any $`0<ϵ<1`$ will do. In this case
$$(\stackrel{~}{f}_k^{(n)}(0))^{\frac{1}{2}}=e^{\frac{i\pi }{n}(\frac{n}{2}+1ϵ)}(f_k^{(n)}(0))^{\frac{1}{2}}=\left|\left(\frac{4}{n}\right)^{\frac{1}{2}}\right|e^{\frac{i\pi }{n}(k\frac{n}{4}ϵ)}.$$
(D.20)
Next we define
$$g_k^{(n)}(z)=F(\stackrel{~}{f}_k^{(n)}(z)),$$
(D.21)
with
$$F(u)=i\frac{1u}{1+u}\frac{au+b}{cu+d},$$
(D.22)
where we use our freedom to fix the signs of $`a,b,c,d`$ to write:
$$\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)=\left(\begin{array}{cc}\frac{1}{\sqrt{2}}e^{\frac{i\pi }{4}}& \frac{1}{\sqrt{2}}e^{\frac{i\pi }{4}}\\ \frac{1}{\sqrt{2}}e^{\frac{i\pi }{4}}& \frac{1}{\sqrt{2}}e^{\frac{i\pi }{4}}\end{array}\right),adbc=1.$$
(D.23)
$`F`$ describes an SL(2,C) map from the disk to the UHP. With this,
$$g_k^{(n)}(0)=\mathrm{tan}\left(\frac{\pi }{n}(k\frac{n}{2}ϵ)\right).$$
(D.24)
For $`k=1,\mathrm{}n`$, the $`g_k^{(n)}(0)`$’s given above are arranged from left to right on the real axis. We also have
$$(g_k^{(n)}(0))^{\frac{1}{2}}=(c\stackrel{~}{f}_k^{(n)}(0)+d)^1(\stackrel{~}{f}_k^{(n)}(0))^{\frac{1}{2}}=\frac{1}{\sqrt{2}}\left|\left(\frac{4}{n}\right)^{\frac{1}{2}}\right|\mathrm{sec}\left(\frac{\pi }{n}(k\frac{n}{2}ϵ)\right).$$
(D.25)
This is manifestly positive for $`1kn`$.
This gives one set of $`g_k^{(n)}`$’s for which the square root rules stated above hold, but we need to show that this holds for any other set of functions $`\stackrel{~}{g}_k^{(n)}(z)`$, related to $`g_k^{(n)}(z)`$ by an SL(2,R) transformation. For this, let us consider another set of functions $`\stackrel{~}{g}_k^{(n)}`$’s related to the $`g_k^{(n)}`$’s via an SL(2,R) transformation $`\left(\begin{array}{cc}p& q\\ r& s\end{array}\right)`$ with the property that $`\stackrel{~}{g}_1^{(n)}(0),\mathrm{}\stackrel{~}{g}_n^{(n)}(0)`$ are arranged from the left to the right on the real axis. In that case, if $`v_k=g_k^{(n)}(0)`$, then for $`k>l`$,
$$v_k>v_l,\frac{pv_k+q}{rv_k+s}\frac{pv_l+q}{rv_l+s}=\frac{(v_kv_l)}{(rv_k+s)(rv_l+s)}>0.$$
(D.26)
Thus
$$(rv_k+s)(rv_l+s)>0.$$
(D.27)
This shows that $`(rv_k+s)`$ has the same sign for all $`k`$. Using the freedom of changing the sign of $`p,q,r,s`$, we can take $`(rv_k+s)`$ to be positive. Then
$$(\stackrel{~}{g}_k^{(n)}(0))^{\frac{1}{2}}=(rv_k+s)^1(g_k^{(n)}(0))^{\frac{1}{2}}>0.$$
(D.28)
This proves the desired result.
Let us now get back to the computation of (D.14) using maps to UHP. For this we map the disk, punctured at $`1,i,1,i`$, into the UHP with the real boundary punctured at $`4,1,0,2`$. These are particularly nice points that give coordinates without radicals:
$`g_1^{(4)}(z)`$ $`=4+6z9z^2+\mathrm{}`$ (D.29)
$`g_2^{(4)}(z)`$ $`=1+\frac{3}{4}z\frac{3}{16}z^2+\mathrm{}`$ (D.30)
$`g_3^{(4)}(z)`$ $`=\mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}0}+\frac{2}{3}z+\frac{1}{9}z^2+\mathrm{}`$ (D.31)
$`g_4^{(4)}(z)`$ $`=\mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}2}+3z+\mathrm{\hspace{0.17em}3}z^2+\mathrm{}`$ (D.32)
In this presentation, all computations are manifestly real. In addition all $`g_i^{(4)}(0)`$’s are positive as expected and we simply take their positive square roots in evaluating (D.14) with $`f_i`$ replaced by $`g_i`$. We get:
$$g^2S|_{t^4}=\frac{(4)(3)+(6)(1)}{123}t^4=\frac{1}{2}t^4$$
(D.33)
This agrees with the result of the disk computation. |
warning/0002/astro-ph0002287.html | ar5iv | text | # Secondary caustics in close multiple lenses
## 1 Introduction
The binary Schwarzschild lens is one of the most intensively studied model. In fact, in a relatively simple way, it shows many features that are observed in general gravitational lenses, such as the formation of multiple images, giant arcs and a not trivial critical behaviour. The first study about the binary lens with equal masses was made by Schneider & Weiß (1986). They derived the critical curves and the caustics showing that three possible topologies are present depending on the distance between the two lenses. Erdl & Schneider (1993) extended these results to a generic mass ratio of the two lenses. Witt & Petters (1993) reached the same results using complex notation. In some limits, Dominik (1999) enlightened the connection between the caustics of the binary lens and other models, such as the Chang–Refsdal lens (Chang & Refsdal 1979; 1984) and the quadrupole lens.
The critical curves and the caustics of multiple lenses can develop very complicated structures, so that the attempts to gain some information about them have been very few. However there is a great interest in this problem for its applications in particular situations, such as planetary systems (Gaudi et al. 1998), rich clusters of galaxies and microlensing of quasars by individual stars in the haloes of the lensing galaxies (Chang & Refsdal 1979; Kayser et al. 1988).
In some special situations, the critical curves of multiple lenses can be derived by perturbative methods, referring to the single Schwarzschild lens as the starting point for series expansions (Bozza 1999; Bozza 2000). These methods work very well in planetary systems, for a lens very far from the others and systems where mutual distances are very small with respect to the total Einstein radius. In the first two cases, the complete caustic structure has been derived and the connections with other models have been showed. In the last case, only the central caustic coming up from the deformation of the total Einstein ring has been studied. Besides this main curve, there are many small critical curves forming among the masses. The caustics generated by these curves generally lie far from the centre of mass and can have some influence on sources distant from the mass distribution. Moreover, they move very quickly as the parameters of the system change (Schramm et al. 1993) constituting the most problematic feature to control in numerical simulations. For these reasons they are sometimes dubbed ghost caustics. In rapidly rotating binaries they may have superluminal projected motion requiring a non–static treatment of light deflection (Zheng & Gould 2000).
In this paper we use complex notation to face the problem of secondary caustics of close multiple lenses. In this way we can study them as deeply as the other caustics, completing the previous works. We shall see that different classes of secondary caustics can be recognized, showing different geometries.
After some review of multiple lensing in Sect. 2, in Sect. 3 we calculate the number and the position of secondary critical curves for an arbitrary number and configuration of lenses. Then, in Sect. 4, we treat the simple caustics and in Sect. 5 the multiple caustics (the distinction will be explained at the end of Sect. 3). In Sect. 6 we specify our formulae for the binary case and in Sect. 7 we give the summary.
## 2 Basics of multiple lensing
We shall study a system of n point–lenses placed at positions $`𝐱_i=(x_{i1};x_{i2})`$ in coordinates normalized to the Einstein radius
$$R_\mathrm{E}^0=\sqrt{\frac{4GM_0}{c^2}\frac{D_{\mathrm{LS}}D_{\mathrm{OL}}}{D_{\mathrm{OS}}}},$$
(1)
where $`M_0`$ is a reference mass (it can be chosen to be the total mass, the typical mass of a single object or anything else). The source coordinates $`𝐲=(y_1;y_2)`$ are normalized to the scaled Einstein radius $`R_\mathrm{E}^0\frac{D_{\mathrm{OS}}}{D_{\mathrm{OL}}}`$. The masses $`m_i`$ of the lenses are measured in terms of $`M_0`$.
We introduce the complex coordinate in the lens plane $`z=x_1+ix_2`$ and the complex source coordinate $`y=y_1+iy_2`$. The positions of the masses will be denoted by $`z_i=x_{i1}+ix_{i2}`$. We also introduce the functions
$$S_k\left(z\right)=\underset{i=1}{\overset{n}{}}\frac{m_i}{\left(zz_i\right)^k}.$$
(2)
The lens equation for our system of n masses reads (Witt 1990)
$$y=z\overline{S}_1\left(\overline{z}\right).$$
(3)
Given a source at position $`y`$, the $`z`$’s solving this equation are the images produced by gravitational lensing.
This map is locally invertible where the determinant of the Jacobian matrix
$$detJ=1\frac{y}{\overline{z}}\overline{\frac{y}{\overline{z}}}=1\left|S_2\left(z\right)\right|^2$$
(4)
is different from zero. The points where the Jacobian determinant vanishes are arranged in smooth closed curves called critical curves. The images of these points through the lens map (3) in the source plane are called caustics. When a source crosses a caustic, creation or destruction of pairs of images occurs and the magnification diverges (Schneider, Ehlers & Falco 1992).
This is all we need to start our search for secondary caustics in close multiple systems. The fundamental hypothesis we make is
$$\left|z_i\right|\sqrt{M}i,$$
(5)
where $`M`$ is the total mass of the system. In this way, the distances between pairs of lenses will be very small with respect to the Einstein radius of the lens that we would have if all the masses were concentrated at the origin. This Einstein radius is $`\sqrt{M}`$ in our notation. The relation (5) allows us to consider the $`z_i`$’s as perturbative parameters in a series expansion. Then we can solve the equation $`detJ=0`$ at each order, writing its solutions as series expansions in powers of the perturbative parameters. In this way we shall find the critical curves of this system and study their properties analytically.
## 3 Number and positions of secondary critical curves
Close multiple lenses have two classes of critical curves: the main critical curve, resulting from the deformation of the Einstein ring of the total mass lens, and the secondary critical curves, forming inside the distribution of the masses.
Effectively, if we multiply the equation $`detJ=0`$ by the quantity $`\underset{i=1}{\overset{n}{}}\left|zz_i\right|^4`$, we get a complex equation in $`z`$ and $`\overline{z}`$:
$$\underset{i=1}{\overset{n}{}}\left|zz_i\right|^4\left|\underset{i=1}{\overset{n}{}}m_i\underset{ji}{}\left(zz_j\right)^2\right|^2=0.$$
(6)
At the zero order, putting all $`z_i`$’s to zero, this equation becomes
$$\left|z\right|^{4n4}\left(\left|z\right|^4M^2\right)=0.$$
(7)
This equation has the solution $`\left|z\right|=\sqrt{M}`$, that is the Einstein ring of the total mass lens. Taking this solution as the starting point of a perturbative expansion, we get the main caustic. The details of this calculation are in (Bozza 2000). But the presence of the solution $`z=0`$ indicates that also this value can be taken as the starting point for another expansion. This is just the value we shall take to find the secondary critical curves.
Having observed the zero order situation, we can start our perturbative approach, searching for the first order solution. Then we write the solution $`z`$ as a series expansion:
$$z=z_0+o\left(\left|z_i\right|\right),$$
(8)
where $`z_0`$ is of the first order in $`\left|z_i\right|`$. Stopping at the first order, we put $`z=z_0`$ in Eq. (6). We see that the first term becomes of order $`4n`$, while the second is of order $`4n4`$. Then the latter dominates the first and Eq. (6) is equivalent to
$$\underset{i=1}{\overset{n}{}}m_i\underset{ji}{}\left(z_0z_j\right)^2=0.$$
(9)
This is a polynomial equation of degree $`2n2`$. Then, for a system of n close lenses, there are, at most, $`2n2`$ points where the Jacobian determinant vanishes (at the first order in $`z_i`$), corresponding to $`2n2`$ secondary critical curves. This is the first main result of our work. It is consistent with the binary lens, since two secondary critical curves are predicted by this formula.
Eq. (9) can be solved analytically for two and three lenses, otherwise we have to resort to simple numerical methods. In Sect. 6, we shall specify these and the following results for the binary lens where a manageable expression for the positions of the critical curves is available. For the triple lens, the analytical solutions are too cumbersome to allow a detailed study.
Now, we have a straightforward way to calculate the positions of the secondary critical curves for an arbitrary configuration of close multiple lenses. Then, we can avoid the traditional blind sampling of the Jacobian determinant on the lens plane and reach, by this new method, the full efficiency.
We take the generical solution $`z_0`$ of Eq. (9) as the first order term of our expansion. From now on, we use the notation
$$S_k^0=S_k\left(z_0\right).$$
(10)
As both $`z_0`$ and $`z_i`$ are of the first order, according to our perturbative expansion, $`S_k^0`$ has all denominators of order $`k`$ and then it is of order $`k`$.
To continue our study we do not need an analytical expression for $`z_0`$. We shall just use the fact that $`z_0`$ is a solution of Eq. (9), that is equivalent to say that
$$S_2^0=0.$$
(11)
Of course, we have to distinguish between simple roots of Eq. (9) and roots of higher multiplicity. Remembering that the $`k^{\mathrm{th}}`$ derivative of $`S_2\left(z\right)`$ is proportional to $`S_{k+2}\left(z\right)`$, we have the equivalence between the following statements:
$$z_0\text{ is a root of multiplicity }pS_{k+2}^0=0k<p.$$
(12)
We shall treat separately the caustics coming from simple roots (hereafter called simple caustics) and the caustics coming from multiple roots (hereafter multiple caustics).
## 4 Simple caustics
These caustics are largely the most common as we explain in the next section. So they surely have the most practical interest.
### 4.1 Shape of the critical curves
Once found the positions of the critical curves, we can carry on our perturbative expansion to discover the shape of these curves. So we put
$$z=z_0+ϵ_2+ϵ_3+\mathrm{},$$
(13)
where $`z_0`$ is the position of one simple critical curve, found by Eq. (9), and $`ϵ_j`$ are the corrections of order $`\left|z_i\right|^j`$. It is convenient to use the original equation $`detJ=0`$, which can be written in the form
$$1S_2\left(z_0+ϵ_2+ϵ_3\right)\overline{S}_2\left(\overline{z}_0+\overline{ϵ}_2+\overline{ϵ}_3\right)=0,$$
(14)
starting from Eq. (4)
The expansion of $`S_2`$ is
$`S_2\left(z_0+ϵ_2+ϵ_3\right)=`$ $`S_2^0+`$ (15)
$`2ϵ_2S_3^0+`$
$`2ϵ_3S_3^0+3ϵ_2^2S_4^0+\mathrm{}.`$
The first row is the order $`2`$ and is null according to Eq. (11). The second row is the order $`1`$ and the third row is the order zero. Inserting this expansion in (14), the lowest order equation is of order $`2`$:
$$4\left|ϵ_2\right|^2\left|S_3^0\right|^2=0,$$
(16)
Being $`z_0`$ a simple root, $`S_3^00`$, so that $`ϵ_2=0`$.
The successive terms in the expansion of the equation (14) are of order zero:
$$14\left|ϵ_3\right|^2\left|S_3^0\right|^2=0.$$
(17)
From this equation we have
$$\left|ϵ_3\right|=\frac{1}{2\left|S_3^0\right|}.$$
(18)
Then the third order contains the first information on the shape of the critical curve. Eq. (18) tells us that the critical curve at position $`z_0`$ is a circle centered on $`z_0`$ with radius
$$r=\frac{1}{2\left|S_3^0\right|}.$$
(19)
From the form of $`S_3^0`$ (see Eq. (2)), we see that the closer the critical curve is to some mass, the higher the value of $`S_3^0`$, the smaller the radius of the circle.
If we multiply all masses by a factor $`\lambda `$, the positions of the critical curves do not change, because $`\lambda `$ factors out from Eq. (9), but their radii change as $`\lambda ^1`$. If we do the same with the positions of the masses instead, the positions of the critical curves scale as $`\lambda `$ and their radii scale as $`\lambda ^3`$.
### 4.2 Caustics
To find the caustics corresponding to the simple critical curves, we just have to put the critical curve, in its obvious parameterization
$$z\left(\theta \right)=z_0+re^{i\theta }\mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}0}\theta <2\pi ,$$
(20)
into the lens equation (3) and expand to the third order:
$$y\left(\theta \right)=\overline{S}_1^0+z_0+\left(re^{i\theta }\frac{e^{2i\theta }}{S_3^0}\right).$$
(21)
We can observe that the lowest order is $`1`$ and is independent on $`\theta `$. It represents the position of the caustic. From the order of this term, we can deduce that these caustics can lie very far from the origin of our system, going to infinity as the distances among the masses are reduced to zero. The successive term is $`z_0`$, which is of the first order and represents a correction to the position. Finally, the shape of the caustic is given by the third order.
The cusps of a caustic are characterized by the vanishing of the tangent vector. To find them, we have to require that
$$\frac{\mathrm{d}y\left(\theta \right)}{\mathrm{d}\theta }=0$$
(22)
and solve for $`\theta `$. Taking $`y\left(\theta \right)`$ from Eq. (21), this equation can be simplified into
$$e^{3i\theta }+\sqrt{\frac{\overline{S}_3^0}{S_3^0}}=0,$$
(23)
whose solutions are
$$\theta _k=\frac{1}{3}\mathrm{arg}\left(S_3^0\right)+\frac{2k\pi }{3}k=0,1,2,$$
(24)
where $`\mathrm{arg}`$ yields the argument of a complex number.
We have three cusps. So, in any close multiple system, having only simple secondary caustics, these caustics have a triangular shape.
Finally, we calculate the area of these caustics. This can be done by the integral
$$A=\underset{\gamma }{}y_2dy_1,$$
(25)
where $`\gamma `$ is the caustic in its clockwise direction. We have
$$A=\frac{1}{4i}\underset{0}{\overset{2\pi }{}}\left[y\left(\theta \right)\overline{y}\left(\theta \right)\right]_\theta \left[y\left(\theta \right)+\overline{y}\left(\theta \right)\right]\mathrm{d}\theta .$$
(26)
The minus in the right member comes from the fact that our parameterization is counterclockwise. The integral only involves complex exponential functions and the result is
$$A=\frac{1}{2}\pi r^2.$$
(27)
So the extension of the simple caustics is of the sixth order in the separations among the lenses, justifying the evasive nature of these caustics.
With our expansions, we have attained considerable analytical information about the secondary caustics establishing their shape, the area, the number of cusps in a completely general way. However, since these results are the fruit of perturbative approximations, it is important to discuss their accuracy. So we propose a comparison between our perturbative results and the numerical ones in a typical situation. We consider a system constituted by three lenses disposed as in Fig. 1. According to our previous statement, this system can form, at most, four simple secondary critical curves. For our choice of parameters, we display their positions in the same figure. The caustics produced by these curves are shown in Fig. 2 where they are compared to the numerical ones. We have taken the distances among the masses of this distribution to be one tenth of the total Einstein radius. Even for this not too small value, the positions and the shapes of the secondary caustics are reproduced with a striking accuracy. It is also to be noted that the quality of numerical results is improved thanks to the guide provided by perturbative results.
So we see that the analytical formulae derived in this section are very good approximations to the quantitative characteristics of the secondary caustics, proving to be highly reliable.
## 5 Multiple caustics
In this section, we consider the case where $`z_0`$ is a multiple root of Eq. (9). The parameters space of a system with n lenses is $`3n4`$ dimensional, since each mass adds three parameters (its mass and its coordinates in the lens plane). Four parameters can be eliminated by considering equivalent those systems differing by a global translation or rotation and/or by a global scale factor. Thus, for example, the binary lens is completely characterized by the mass ratio and the separation between the lenses.
The requirement of a double root in Eq. (9) translates into the vanishing of the derivative of this equation with respect to $`z`$. This is one constraint equation, then the points of the parameters space producing multiple roots constitute a $`3n5`$ dimensional hypersurface, thus having measure zero. For this reason, the occurrence of multiple roots is relatively rare. Anyway, very interesting features emerge, justifying a detailed study of these particular cases.
### 5.1 Critical curves
Suppose that $`z_0`$ is a root with multiplicity $`p`$. We have to find the correct order of the perturbation to insert in the equation $`detJ=0`$, representing the shape of our critical curve. According to the equivalence (12), the $`S_{k+2}^0`$ with $`k<p`$ are null. Then, we put
$$z=z_0+ϵ,$$
(28)
with the order (that we shall indicate by $`q`$) of $`ϵ`$ to be found. We only assume that $`q`$ be higher than one. Then, the expansion of $`S_2\left(z\right)`$ is
$$\begin{array}{c}S_2\left(z_0+ϵ\right)=S_2^02ϵS_3^0+3ϵS_4^0+\mathrm{}\hfill \\ \hfill \mathrm{}+\left(1\right)^k\left(k+1\right)ϵ^kS_{k+2}^0+\mathrm{}.\end{array}$$
(29)
The $`k^{\mathrm{th}}`$ term is of order $`qk\left(k+2\right)`$ and the first term to be non–null is that for $`k=p`$. When we put this expansion in the equation $`detJ=0`$, the first non–null term is
$$\left(p+1\right)^2\left|ϵ^pS_{p+2}\right|^2$$
(30)
having order $`2qp2\left(p+2\right)`$. If this order is less than zero, we just get from $`detJ=0`$ that $`ϵ=0`$, but if the order of this term is zero, then the zero order expansion of $`detJ=0`$ also involves another term (equal to $`1`$):
$$1\left(p+1\right)^2\left|ϵ^pS_{p+2}\right|^2=0$$
(31)
and the equation gives the non–trivial solution
$$\left|ϵ\right|=\frac{1}{\left[\left(p+1\right)\left|S_{p+2}\right|\right]^{1/p}}.$$
(32)
This happens when the order of $`ϵ`$ is $`q=\frac{2+p}{p}`$. This is consistent with the result of the previous section, because, for $`p=1`$, $`q=3`$. For $`p=2`$, we have that the first non trivial order is the second and, for $`p=3`$, it is the order $`\frac{5}{3}`$. When $`p`$ increases, the order of this perturbation decreases, approaching 1 as a limit. This means that at high multiplicities, the perturbative expansion becomes always less accurate, requiring ever more terms for an adequate description of the caustics. Anyway, the main characteristics of the caustics can be derived retaining just the first correction and that is what we shall do.
The critical curve just derived is again a circle with radius
$$r=\frac{1}{\left[\left(p+1\right)\left|S_{p+2}\right|\right]^{1/p}},$$
(33)
becoming greater with the multiplicity.
### 5.2 Caustics
We take, as before, the parameterization
$$z\left(\theta \right)=z_0+re^{i\theta }$$
(34)
for the critical curve, with $`r`$ given by Eq. (33). Putting this expression into the lens equation (3) and expanding to the $`q^{\mathrm{th}}`$ order, we get
$$y\left(\theta \right)=\overline{S}_1^0+z_0+r\left[e^{i\theta }+\left(1\right)^p\frac{e^{\left(p+1\right)i\theta }}{p+1}\sqrt{\frac{\overline{S}_{p+2}^0}{S_{p+2}^0}}\right].$$
(35)
The $`q^{\mathrm{th}}`$ order is the first depending on $`\theta `$ and determines the shape of the caustic.
To understand this shape, we calculate the cusps as in the previous section. The equation for the cusps is
$$e^{\left(p+2\right)i\theta }+\left(1\right)^{p+1}\sqrt{\frac{\overline{S}_{p+2}^0}{S_{p+2}^0}}=0$$
(36)
and its solutions are
$$\theta _k=\frac{\left(1\right)^p}{p+2}\mathrm{arg}\left(S_{p+2}^0\right)+\frac{2k\pi }{p+2}\mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}0}k<p+2.$$
(37)
Now we have $`p+2`$ cusps. This is a very interesting result, because the caustic assumes the shape of a regular polygon with $`p+2`$ curved sides.
The area of the multiple caustic can be calculated in the same way as for the simple one. We just give the result:
$$A=\frac{p}{p+1}\pi r^2.$$
(38)
It is of order $`2\frac{2+p}{p}`$. So, increasing the multiplicity from 1 to infinity, the order of the area lowers from 6 to 2 and the extension of the caustic becomes ever more important.
Some other consideration about the limit for $`p\mathrm{}`$ can be done. The number of cusps become infinite and, from Eq. (35), we see that the caustic becomes a circle of radius $`r`$. In fact, the area becomes $`\pi r^2`$.
### 5.3 An example: a double caustic in a triple lens
Now we shall practically see how our formulae work in the case of a multiple caustic. We consider three masses: $`m_1=0.25`$, $`m_2=0.25`$ and $`m_3=0.5`$. We fix the positions of the first two: $`z_1=0.1`$, $`z_2=0.1`$; but we let the third free for the moment. We simultaneously solve Eq. (9) and its derivative with respect to $`z_0`$, for the two unknowns $`z_0`$ and $`z_3`$. We find six possible positions of the third mass, giving rise to a double root of the positions equation. None of them is a triple root. Two of these positions are on the $`x_2`$-axis. We choose one of them: $`z_3=i0.084263`$. The double root is in $`z_0=i0.0299`$. In Fig. 3, we see that the critical curve in this point is much greater than the other two. In fact, the radius of the double critical curve, calculated by Eq. (33), is $`8.49\times 10^3`$, while the radius of the two simple critical curves is $`6.13\times 10^4`$, according to Eq. (19).
In Fig. 4, we show the caustic generated by this double critical curve. The geometry is correctly predicted by our perturbative expansion: there are four cusps in a double caustic. We see that the approximation is less accurate than before, as we anticipated in our discussion about the order of the perturbation. However, for double caustics, it is not so difficult to add another term to the perturbative expansion and reach the same accuracy of the simple caustics. The third order term in the critical curve depends on $`\theta `$:
$$ϵ_3=6r^2\mathrm{Re}\left[\overline{S}_4^0S_5^0e^{i\theta }\right].$$
(39)
The successive term in the caustic is
$$e^{i\theta }ϵ_3+3e^{3i\theta }r^2ϵ_3\overline{S}_4^0r^4e^{4i\theta }\overline{S}_5^0.$$
(40)
Double caustics, and, more generally, multiple caustics, are formed by the union of small caustics, in some sense. Another interesting question is: what happens if we change the parameters in the neighbourhood of our particular choice producing the double root? We expect the double critical curve to separate into two smaller ovals and the double quadrangular caustic to break into two triangular ones; but this can happen in different ways.
In this regime, the perturbative caustics are simple. However, as the parameters tend to give the double root, $`S_3^0`$ tends to zero, yielding a diverging $`r`$ for the simple critical curves, according to Eq. (19). The transition with the formation of the double critical curve is thus not reproduced. Guided by perturbative approximations, the break of the double caustic, when $`z_3`$ moves out from the position $`i0.084263`$, can be investigated numerically. The results are shown in Fig. 5.
In case $`a`$, $`z_3=i0.082787`$, i.e. we have moved the third mass towards the others. The critical curve breaks in the horizontal direction. Looking just at the thick line in Fig. 5a2, representing the numerical caustic, we see that the top cusp and the bottom cusp develop a butterfly geometry. At some critical value, these butterflies touch and the two resulting triangular caustics move away along the horizontal direction. We have displayed in the same plot the perturbative caustics too. Obviously, they are simple caustics, so they cannot show the butterfly geometry but they can help in understanding how the separation occurs. We also notice that the simple caustics cover the area of the numerical transition double caustic very well constituting a significant approximation anyway.
In case $`b`$, $`z_3=i0.085759`$, so that the third mass is farther from the others. Now the critical curve breaks in the vertical direction and so does the caustic. The left and the right cusps transform into butterflies. These butterflies are slightly distorted by the fact that the resulting simple caustics have different sizes: the one on the top is smaller than the other.
In case $`c`$, $`z_3=0.00147+i0.084263`$. We have displaced the third mass in the horizontal direction. The critical curve breaks diagonally and so does the caustic. But this time the transition occurs with a simple beak–to–beak singularity rather than with butterflies. While in the previous situations the two simple caustics in the last step of the separation touch with a fold, here they touch with a cusp.
## 6 Secondary caustics in binary lensing
In this section, we specify our results for the binary case where simple analytical formulae can be written. Let’s consider two masses placed on the horizontal axis and let’s choose the origin in the centre of mass. We call the separation between the masses $`a`$, then we have $`z_1=\frac{m_2a}{m_1+m_2}`$ and $`z_2=\frac{m_1a}{m_1+m_2}`$. Eq. (9) is of second degree. Its solutions are
$$z_0=a\frac{m_2m_1\pm i\sqrt{m_1m_2}}{m_1+m_2}.$$
(41)
They are always simple and lie on a circle of radius $`a/2`$ centered in the middle of the two masses.
The radius of the two ovals is the same:
$$r=\frac{\sqrt{m_1m_2}a^3}{2\left(m_1+m_2\right)^2}.$$
(42)
Its maximum value $`\frac{a^3}{4M_{\mathrm{tot}}}`$ is reached when the two masses are equal, in fact, in this case, their distance from the two masses is maximum.
The two caustics are given by the following expression:
$$\begin{array}{c}y\left(\theta \right)=\frac{m_1m_2}{a}i\frac{2\sqrt{m_1m_2}}{a}+\hfill \\ \hfill +a\frac{m_2m_1}{m_1+m_2}\pm ia\frac{\sqrt{m_1m_2}}{m_1+m_2}+\\ \hfill +a^3\sqrt{m_1m_2}\left[\frac{e^{i\theta }}{2\left(m_1+m_2\right)^2}\pm \frac{ie^{2i\theta }}{4\left(\sqrt{m_1}\pm i\sqrt{m_2}\right)^4}\right].\end{array}$$
(43)
Their cusps are at positions
$$\theta _k=\mathrm{arg}\left[\pm i\left(\sqrt{m_1}\pm i\sqrt{m_2}\right)^4\right]+\frac{2k\pi }{3}k=0,1,2.$$
(44)
Their area is
$$A=\frac{\pi m_1m_2a^6}{8\left(m_1+m_2\right)^4},$$
(45)
reaching the maximum value $`\frac{\pi a^6}{32M_{\mathrm{tot}}^2}`$ in the equal masses case.
## 7 Summary
Multiple Schwarzschild lensing represents a very rich terrain for the exploration of caustics in gravitational lensing. The occurrence of different kinds of singularities stimulates new investigations.
In this paper we have applied perturbative methods to secondary caustics, forming when the masses are close each other with respect to the total Einstein radius. In this way we have been able to establish the number of the caustics for any lens configuration, the positions and the shapes, with a complete characterization of the geometries arising in all cases. Moreover, quantitative formulae for the area and other features of these objects have been given. As we have seen, in the most common case, the shape of the simple caustics is always triangular. Anyway, multiple caustics exist developing a great variety of behaviours, giving rise, to curves having a number of cusps ranging from four to infinity. The breaking of multiple caustics can follow different ways depending on how the parameters of the system change.
###### Acknowledgements.
I would like to thank Gaetano Scarpetta and Salvatore Capozziello for their helpful comments on the manuscript. Work supported by fund ex 60% D.P.R. 382/80. |
warning/0002/astro-ph0002010.html | ar5iv | text | # RADIATION SPECTRA FROM ADVECTION-DOMINATED ACCRETION FLOWS IN A GLOBAL MAGNETIC FIELD
## 1. INTRODUCTION
<sup>1</sup><sup>1</sup>footnotetext: Present adress: Department of Earth and Space Science, Graduate School of Science, Osaka University, Toyonaka, Osaka 560-0043, Japan
kino@vega.ess.sci.osaka-u.ac.jp
The optically thin, advection-dominated accretion flows (ADAFs) have been studied by a number of authors during past several years (e.g. Narayan $`\&`$ Yi 1994, 1995a, b; Abramowicz et al. 1995; Nakamura et al. 1997, Manmoto, Mineshige & Kusunose 1997; Narayan et al. 1998). These models are very successful in describing both spectra and dynamics of accreting black hole systems such as those in binaries and in low-luminosity active galactic nuclei (AGNs). The observed spectra can be explained as follows. The radio emission is due to the synchrotron emission in turbulent magnetic fields in the accretion flow. These synchrotron photons serve as seed photons for the inverse Compton process by hot electrons. Once-scattered Compton photons are mainly distributed in the optical band and twice-scattered Compton photons, in soft X-ray band. Bremsstrahlung due to electron-electron and electron-proton collisions gives rise to the observed hard X-ray spectra.
Thus, these ADAF models provide a good framework for understanding the observed spectra. In these models, both angular momentum transfer and energy dissipation in the accretion flow is assumed to be undertaken by the turbulent viscosity whose size is specified by so-called $`\alpha `$ parameter. For this reason, hereafter we call this type of models the “viscous” ADAF model in this paper. The magnetic fields are regarded as of turbulence origin and are described by another parameter $`\beta `$ which specifies the ratio of the magnetic pressure to the gas pressure.
However there is no reason to believe that the turbulent viscosity is the only candidate that controls the accretion processes. Rather, it is quite natural to think that some types of global magnetic fields may play an essential role. Indeed, there are some evidences for the presence of such an ordered magnetic field in the central region of our Galaxy (e.g., Yusef-Zadeh, Morris & Chance 1984). As Kato, Fukue & Mineshige (1998) has pointed out the hydromagnetic turbulence in accretion disks may also generate global magnetic fields by dynamo processes due to the presence of helical motions. In view of such circumstances, another type of ADAF model has been proposed by one of the present authors (Kaburaki 1999, 2000; hereafter referred to as K99 and K00). In order to distinguish it from the above viscous ADAF models, hereafter we call it the “resistive” ADAF models since energy dissipation in the accretion flow is due to the electric resistivity and angular momentum transfer is supported not by the viscosity but by the magnetic stress of a large scale magnetic field.
The purpose of the present study is to calculate the expected radiation spectra from ADAFs in a global magnetic field based on the resistive ADAF model, in order to compare its predictions with those of the viscous ADAF models. As a most suitable candidate for such a comparison, Sgr A is taken up here because it has been observed in many wave lengths as the nearest galactic nucleus and its spectrum has been reproduced many times by the successively advancing viscous ADAF models.
In §2, we introduce the set of analytic solutions for resistive ADAFs in a suitably scaled form and discuss their basic characteristics. The relevant radiation mechanisms and the methods of calculation of the fluxes are described in §3. These schemes are applied to Sgr A in §4 and the results are discussed in comparison with those of the viscous ADAF models. Finally in §5, we summerize the main results and discuss some related issues.
## 2. RESISTIVE ADAF SOLUTION
As a basis of our calculation of spectra, we introduce here the set of analytic solutions constructing the resistive ADAF model. This set may be considered as a counterpart of that found by Narayan $`\&`$ Yi (1994, 1995a) in the viscous ADAF scheme, but it should be emphasized that the former is not a self-similar solution as the latter. In the resistive ADAF model, there are three basic quantities and one parameter: mass of the central black hole $`M`$, mass accretion rate $`\dot{M}`$, strength of the external magnetic field $`|B_0|`$ and half-opening angle of the flow $`\mathrm{\Delta }`$, respectively. We introduce the following normalizations for these quantities and for the radial distance $`R`$:
$$m\frac{M}{10^6M_{}},\dot{m}\frac{\dot{M}}{\dot{M}_\mathrm{E}},b_0\frac{|B_0|}{1G},\delta \frac{\mathrm{}}{0.1},r\frac{R}{R_{\mathrm{out}}},$$
(1)
where $`R_{\mathrm{out}}`$ denotes the radius of the disk’s outer edge. The Eddington accretion rate is defined by $`\dot{M}_\mathrm{E}L_\mathrm{E}/(0.1c^2)`$, which includes the efficiency factor of $`0.1`$. Note that this definition of $`\dot{M}_\mathrm{E}`$ and the normalization factor for black hole mass are different from those in K00. The latter is chosen here as 10<sup>6</sup> $`M_{}`$ for the convenience of the discussion of Sgr A.
In spherical polar coordinates, the radius of outer edge and the radial-part functions of relevant physical quantities (for their angular parts, see also K99, K00) are written as
$$R_{\mathrm{out}}=1.5\times 10^{16}b_0^{4/5}\dot{m}^{2/5}m^{3/5}\mathrm{cm},$$
(2)
$$|b_\phi (r)|=10\delta ^1b_0r^1\mathrm{G},$$
(3)
$$v_\mathrm{K}(r)=0.9\times 10^8b_0^{2/5}\dot{m}^{1/5}m^{1/5}r^{1/2}\mathrm{cm}\mathrm{s}^1,$$
(4)
$$P(r)=4.0\delta ^2b_0^2r^2\mathrm{dyne}\mathrm{cm}^2,$$
(5)
$$T(r)=1.8\times 10^7b_0^{4/5}\dot{m}^{2/5}m^{2/5}r^1\mathrm{K},$$
(6)
$$\rho (r)=1.3\times 10^{15}\delta ^2b_0^{6/5}\dot{m}^{2/5}m^{2/5}r^1\mathrm{g}\mathrm{cm}^3,$$
(7)
where $`b_\phi `$ is the toroidal magnetic field, $`v_\mathrm{K}`$ is the Kepler velocity, $`P`$ is the gas pressure, $`T`$ is the temperature (common to electrons and ions) and $`\rho `$ is the density. Owing to the non-negligible pressure term in the radial force balance, the toroidal velocity in the disk is reduced by a factor of $`1/\sqrt{3}`$ from the Kepler value. The surface density and the optical depth are given, respectively, by
$$\mathrm{\Sigma }(r)=\mathrm{\Sigma }=4.1\delta ^1b_0^{2/5}\dot{m}^{4/5}m^{1/5}\mathrm{g}\mathrm{cm}^2,$$
(8)
$$\tau _{\mathrm{es}}(r)\frac{1}{2}\kappa _{\mathrm{es}}\mathrm{\Sigma }=8.2\times 10^1\delta ^1b_0^{2/5}\dot{m}^{4/5}m^{1/5},$$
(9)
where $`\kappa _{\mathrm{es}}`$ is the opacity for electron scattering. Note that these are independent of $`r`$ in the present model.
Fig. 1 shows a schematic picture of an accretion disk in a global magnetic field, whose precise structure is described by the solution given above. Otherwise uniform external magnetic field is twisted by the rotational motion of accreting plasma, and there develops a large toroidal magnetic field in the middle latitude region. The behavior of this component, especially within the geometrically thin accretion flow, is given as $`b_\phi \mathrm{tanh}\xi `$ in the resistive ADAF solution, where $`\xi =(\theta \pi /2)/\mathrm{\Delta }`$ is the normalized angular variable. Owing to the appearance of a global $`b_\phi `$, angular momentum of the accreting plasma becomes able to be carried away by the magnetic stress to distant regions along the poloidal magnetic field. The extraction of angular momentum guarantees the inward motion of the plasma, which gradually becomes large until it reaches near the rotational velocity at around the inner edge of the accretion disk. Although the magnetic lines of force are also bent inwardly toward the center of gravitational attraction, it has been shown that the dominant component is the toroidal one. This component also plays an essential role in the plasma confinement toward the equatorial plane and keeps it geometrically thin through its magnetic pressure.
Before going into the detailed discussion of the radiative processes, we briefly mention some similarities and differences in the basic features of the viscous and resistive ADAF models. It is worth noting that, in spite of the essential difference in the mechanisms of angular momentum transport and energy dissipation, the predictions for quantities such as temperature, density and optical depth are quite similar in both models. The temperatures in both models are as high as a fraction of the virial temperature of ions. Indeed, this is the case in the viscous models, though the electron temperature may deviates from it in the inner regions (e.g., Narayan & Yi 1995b), and so also in the resistive model as can be confirmed from the analytic expression of $`T`$ (K99, K00). Such a high temperature makes a sharp contrast with the case of standard $`\alpha `$-disks (e.g., Frank, King $`\&`$ Raine 1992). Further, for a sub-Eddington mass accretion rate, the optical depth is dominated by the electron scattering and is smaller than unity. These are therefore common features of sub-Eddington ADAF models as expected.
As one of the main differences, it may be stressed that the magnetic field in the resistive ADAF model is an ordered magnetic field and is determined self-consistently in the model from a boundary value. Therefore, the strength of the field is not a parameter as in the viscous ADAF models. The ordered magnetic filed extract angular momentum from the accreting plasma and confines it in a disk structure against the gas pressure. Gravitational energy is released in the disk as the Joule heating and also as compressional heating of the flow. In the above analytic model of resistive ADAFs, these energies are fully advected down the stream. We calculate the radiation from the disk as a small perturbation from this solution.
Another distinction may be in the energy partition between the electron and ion components. The viscous ADAF models assume that the viscous dissipation, which is large at large radii, heats mainly ions. Since the efficiency of radiation cooling is very small for ions compared with electrons and since energy transfer to the electron component is estimated to be negligible (Manmoto et al. 1997) except in outer portions, the flow is fully advective in most portions. For the electron component, on the other hand, the radiative cooling is balanced by the advective heating, near the inner edge. The electron temperature, therefore, deviates downwards largely from the ion temperature thus realizing a two-temperature structure.
In contrast to the viscous heating, the resistive dissipation becomes large at small radii and seems to preferentially heat the electron component as suggested by Bisnovatyi-Kogan & Lovelace (1997). In this case, the temperature difference is expected to remain rather small because the heating is effective for effective radiator. In any case, the resistive ADAF model in its present version assume a common temperature to both components, for simplicity. The examination of its two-temperature version may belong to a future work.
## 3. CALCULATION OF SPECTRUM
As mentioned above, the radiation spectrum from a resistive ADAF is calculated based on the analytic solution introduced in the previous section. Back reactions of the radiation cooling to this fully advective solution are negligible as far as its fraction in the total cooling rate is small. This has been roughly checked in a previous paper (K00). The method of calculating radiation fluxes described in this section will be applied to Sgr A in the next section. The observed spectrum of Sgr A in the frequency range from radio up to X-ray range is successfully explained in the viscous ADAF models by the three processes, i.e., synchrotron radiation, bremsstrahlung and inverse Compton scattering (Narayan, Yi & Mahadevan 1995; Manmoto et al. 1997; Narayan et al. 1998). Although there may be some other components such as the radio-frequency excess over the Rayleigh-Jeans spectrum and $`\gamma `$-ray peak both of which need separate explanations (see, e.g., Mahadevan 1999 for the former, and Mahadevan, Narayan & Krolik 1997 for the latter), we ignore these components here for simplicity.
Among the above three processes, the Compton scattering is treated separately from the other processes of emission and absorption. Therefore, we divide the total flux into two parts: the flux due to bremsstrahlung and synchrotron process $`F_\nu `$ and that due to the inverse Compton process $`F_\nu ^{}`$. The obtained fluxes are both integrated over the entire surfaces (upper and lower ones) of a disk, and added up to obtain the luminosity per unit frequency $`L_\nu `$.
Temperature in the flow is vertically isothermal in the present model. In calculating the emission and absorption processes, the flow is further assumed to be locally plane parallel. Solving the radiative diffusion equation at a given radius $`R`$, we obtain the flux of the unscattered photons $`F_\nu `$ emanating from one side of the disk (Rybicki & Lightman 1979) as
$$F_\nu =\frac{2\pi }{\sqrt{3}}B_\nu \left[1\mathrm{exp}(2\sqrt{3}\tau _\nu ^{})\right],$$
(10)
where $`B_\nu `$ is the Planck intensity and $`\tau _\nu ^{}`$ is the vertical optical depth for absorption,
$$\tau _\nu ^{}(R)\frac{\sqrt{\pi }}{2}\kappa _\nu R\mathrm{\Delta }.$$
(11)
Assuming the local thermodynamic equilibrium (LTE), we can express the absorption coefficient $`\kappa _\nu `$ at the equatorial plane in terms of the volume emissivities $`\chi _\nu `$’s for bremsstrahlung and synchrotron processes:
$$\kappa _\nu =\frac{\chi _{\nu ,\mathrm{br}}+\chi _{\nu ,\mathrm{sy}}}{4\pi B_\nu }.$$
(12)
Thus, equation (10) includes not only the effect of free-free absorption but also of synchrotron self-absorption at low frequencies.
As the distribution function for thermal electrons, we assume that of the relativistic Maxwellian (in its normalized form),
$$N_e(\gamma )d\gamma =\frac{\gamma ^2\beta \mathrm{exp}(\frac{\gamma }{\theta _e})}{\theta _eK_2(\frac{1}{\theta _e})}d\gamma ,\theta _e=\frac{k_\mathrm{B}T_e}{m_ec^2},$$
(13)
because ADAFs tend to have so high temperatures that electron thermal energy can exceed its rest mass energy. Here, $`\gamma `$ is the Lorentz factor, $`k_\mathrm{B}`$ is the Boltzmann constant and $`K_2`$ is the 2nd modified Bessel function. Actually, we use this formula only in the calculation of Comptonized photon flux below, while in those of bremsstrahlung and synchrotron processes we follow the works of Narayan & Yi (1995b), and Manmoto et al. (1997) where it is replaced by a numerical fitting function.
### 3.1. Bremsstrahlung
At relativistic temperatures, we must take into account not only electron-proton but also electron-electron encounters. Therefore, the total bremsstrahlung cooling rate per unit volume is written as
$$q_{\mathrm{br}}^{}=q_{ei}^{}+q_{ee}^{},$$
(14)
where the subscripts $`ei`$ and $`ee`$ denote the electron-ion and electron-electron processes, respectively. The explicit expressions of the cooling rates are as follows. For the electron-proton process,
$`q_{ei}^{}`$ $`=`$ $`1.25n_e^2\sigma _\mathrm{T}c\alpha _fm_ec^2F_{ei}(\theta _e)`$ (15)
$`=`$ $`1.48\times 10^{22}n_e^2F_{ei}(\theta _e)\mathrm{ergs}\mathrm{cm}^3\mathrm{s}^1,`$
where $`n_e`$ is the electron number density, $`\alpha _f`$ is the fine-structure constant and $`\sigma _\mathrm{T}`$ is the Thomson cross-section, and further
$`F_{ei}(\theta _e)`$ $`=`$ $`4\left({\displaystyle \frac{2\theta _e}{\pi ^3}}\right)^{0.5}(1+1.781\theta _e^{1.34})\mathrm{for}\theta _e<1,`$ (16)
$`=`$ $`{\displaystyle \frac{9\theta _e}{2\pi }}[\mathrm{ln}(1.123\theta _e+0.48)+1.5]`$
$`\mathrm{for}\theta _e>1.`$
For the electron-electron process,
$`q_{ee}^{}`$ $`=`$ $`n_e^2cr_e^2m_ec^2\alpha _f{\displaystyle \frac{20}{9\pi ^{0.5}}}(443\pi ^2)\theta _e^{3/2}`$ (17)
$`\times (1+1.1\theta _e+\theta _e^21.25\theta _e^{5/2})`$
$`=`$ $`2.56\times 10^{22}n_e^2\theta _e^{3/2}(1+1.1\theta _e+\theta _e^21.25\theta _e^{5/2})`$
$`\mathrm{ergs}\mathrm{cm}^3\mathrm{s}^1`$
when $`\theta _e<1`$, and
$`q_{ee}^{}`$ $`=`$ $`n_e^2cr_e^2m_ec^2\alpha _f24\theta _e(\mathrm{ln}2\eta \theta _e+1.28)`$ (18)
$`=`$ $`3.40\times 10^{22}n_e^2\theta _e(\mathrm{ln}1.123\theta _e+1.28)`$
$`\mathrm{ergs}\mathrm{cm}^3\mathrm{s}^1`$
when $`\theta _e>1`$. Here, $`r_e=e^2/m_ec^2`$ is the classical electron radius and $`\eta =\mathrm{exp}(\gamma _E)=0.5616`$.
The emissivity per frequency is given by
$$\chi _{\nu ,\mathrm{br}}=q_{\mathrm{br}}^{}\overline{G}\mathrm{exp}\left(\frac{h\nu }{k_\mathrm{B}T_e}\right)\mathrm{ergs}\mathrm{cm}^3\mathrm{s}^1\mathrm{Hz}^1,$$
(19)
where $`h`$ is the Planck constant and $`\overline{G}`$ is the Gaunt factor which is written (Rybicki $`\&`$ Lightman 1979) as
$`\overline{G}`$ $`=`$ $`{\displaystyle \frac{h}{k_\mathrm{B}T_e}}\left({\displaystyle \frac{3}{\pi }}{\displaystyle \frac{k_\mathrm{B}T_e}{h\nu }}\right)^{1/2}\mathrm{for}{\displaystyle \frac{h\nu }{k_\mathrm{B}T_e}}>1,`$ (20)
$`=`$ $`{\displaystyle \frac{h}{k_\mathrm{B}T_e}}{\displaystyle \frac{\sqrt{3}}{\pi }}\mathrm{ln}\left({\displaystyle \frac{4}{\zeta }}{\displaystyle \frac{kT_e}{h\nu }}\right)\mathrm{for}{\displaystyle \frac{h\nu }{k_\mathrm{B}T_e}}<1.`$
The above cited formule contain a few minor defects. For example, the non-relativistic limit calculated for electron-ion process from equations (16) and (20) differs by about 35% from the standard formula (Rybicki & Lightman 1979). Equation (20) assumes the same values of the Gaunt factor for both electron-electron and electron-ion processes. In spite of these defects, we adopt the above formule according to Narayan & Yi (1995b) and Manmoto et al. (1997), considering that these are the best ones we can employ at present throughout the energy rage of our interest. The adoption of the same formule as in the previous calculations is also suitable for the purpose of comparison of the predictions of different models, such as the viscous and resistive ones.
### 3.2. Synchrotron Emission
Synchrotron emission is an essential process to produce the radio wave-length part of the spectra from optically thin ADAFs in AGNs. Especially in the resistive ADAF model, some information about the strength of the ambient magnetic field may be obtained from the process of spectral fitting.
The optically-thin synchrotron emissivity by relativistic Maxwellian electrons is calculated from the formula (Narayan & Yi 1995b; Mahadevan, Narayan & Yi 1996),
$`\chi _{\nu ,\mathrm{sy}}`$ $`=`$ $`4.43\times 10^{30}{\displaystyle \frac{4\pi n_e\nu }{K_2(1/\theta _e)}}I^{}\left({\displaystyle \frac{4\pi m_ec\nu }{3eB\theta _e^2}}\right)`$ (21)
$`\mathrm{ergs}\mathrm{cm}^3\mathrm{s}^1\mathrm{Hz}^1,`$
where $`e`$ is the elementary charge and
$`I^{}(x)={\displaystyle \frac{4.0505}{x^{1/6}}}\left(1+{\displaystyle \frac{0.4}{x^{1/4}}}+{\displaystyle \frac{0.5316}{x^{1/2}}}\right)\mathrm{exp}(1.8899x^{1/3}).`$ (22)
In equation(21), the argument of $`I^{}`$ is specified as
$`x{\displaystyle \frac{2\nu }{3\nu _0\theta _e^2}},\nu _0{\displaystyle \frac{e|B|}{2\pi m_ec}},`$ (23)
where $`B`$ is the local value of magnetic field for which we substitute $`b_\phi `$.
### 3.3. Inverse Compton Scattering
The soft photons whose flux is given by equation (10) are Compton scattered by the relativistic electrons in the flow. We adopt the rate equation of Coppi & Blandford (1990) as the basis of our considerations. This equation applies to homogeneous, isotropic distributions. The first term on the right-hand side of their equation describes the rate of decrease in the photon’s number density with a given energy owing to the scattering into other energies, while the second term does the increase owing to the scattering into this energy from other energies.
In the situations of our interest, we can neglect the first term because the number density of Comptonized photons are small compared with that of the seed photons. Instead, we use the second term iteratively to calculate the effects of multiple scattering. The scattering occurs on the average when the condition $`c\sigma _\mathrm{T}n_edt=1`$ is satisfied, where $`t`$ is time and $`n_e`$ is the number density of electrons. The probability that such a condition is satisfied $`j`$-times before the photons come out of the surface may be given by the Poisson formula,
$`p_j={\displaystyle \frac{\tau _e^j\mathrm{e}^{\tau _e}}{j!}}.`$ (24)
Then, the production rate for the photons with a normalized energy $`ϵh\nu /m_ec^2`$ is given by
$`{\displaystyle \frac{dn(ϵ)}{c\sigma _\mathrm{T}n_edt}}={\displaystyle \underset{j=1}{\overset{\mathrm{}}{}}}p_j{\displaystyle 𝑑\gamma _j\mathrm{}𝑑\gamma _1N_e(\gamma _j)\mathrm{}N_e(\gamma _1)}`$ (25)
$`\times {\displaystyle }dϵ_j\mathrm{}dϵ_1`$
$`[P(ϵ;ϵ_j,\gamma _j)\mathrm{}P(ϵ_2;ϵ_1,\gamma _1)R(ϵ_j,\gamma _j)\mathrm{}`$
$`R(ϵ_1,\gamma _1)n_{\mathrm{in}}(ϵ_1)],`$
where $`m_e`$ is the electron mass and $`n_{\mathrm{in}}`$ is the number density of seed photons.
The non-dimensional scattering rate $`R(ϵ,\gamma )`$ including Klein-Nishina cross section $`\sigma _{\mathrm{KN}}`$ is written explicitly (Coppi $`\&`$ Blandford 1990) as
$`R(ϵ,\gamma )`$ $`=`$ $`{\displaystyle _1^1}{\displaystyle \frac{d\mu }{2}}(1\beta \mu ){\displaystyle \frac{\sigma _{\mathrm{KN}}(\beta ,ϵ,\mu )}{\sigma _\mathrm{T}}}`$ (26)
$`=`$ $`{\displaystyle \frac{3}{32\gamma ^2\beta ϵ^2}}{\displaystyle _{2\gamma (1\beta )ϵ}^{2\gamma (1+\beta )ϵ}}dx[(1{\displaystyle \frac{4}{x}}{\displaystyle \frac{8}{x^2}})\mathrm{ln}(1+x)`$
$`+{\displaystyle \frac{1}{2}}+{\displaystyle \frac{8}{x}}{\displaystyle \frac{1}{2(1+x)^2}}]\mathrm{cm}^3\mathrm{s}^1.`$
Scattered-photon distribution is denoted by $`P(ϵ;ϵ^{},\gamma )`$ and, in the present calculation, approximated by a $`\delta `$-function (Lightman & Zdziarski 1987, Fabian et al. 1986):
$`P(ϵ;ϵ^{},\gamma )=\delta \left(ϵ{\displaystyle \frac{4\gamma ^2}{3}}ϵ^{}\right).`$ (27)
This is merely for simplicity and a more exact expression has been derived by Jones (1968) and corrected afterwards by Coppi & Blandford (1990).
Although the sum in equation (25) runs to infinity, it seems appropriate to assume that photons which are scattered more than certain times become saturated and obey the Wien distribution $`\nu ^3\mathrm{exp}(h\nu /kT_e)`$ (e.g., Manmoto et al. 1997). In view of the smallness of the optical depth in most sub-Eddington ADAFs ($`\tau _{\mathrm{es}}<10^3`$, in the case of Sgr A), however, we truncate the power series in $`\tau _{\mathrm{es}}`$ at $`j=2`$ and ignore the saturation effect. After performing the integrations containing $`\delta `$-functions and transforming the photon number densities into fluxes by multiplying $`chϵ/2`$ on both sides of equation (25), we obtain
$$F_\nu ^{}(0)=\mathrm{e}^{\tau _{\mathrm{es}}}\left[\tau _{\mathrm{es}}F_\nu ^{(1)}+\frac{\tau _{\mathrm{es}}^{2}}{2}F_\nu ^{(2)}\right],$$
(28)
where once- and twice-scattered fluxes are given, respectively, by
$$F_\nu ^{(1)}_1^{\mathrm{}}𝑑\gamma _1N_e(\gamma _1)R(\frac{3ϵ}{4\gamma _1^{2}},\gamma _1)F_{\mathrm{in}}\left(\frac{3ϵ}{4\gamma _1^{2}}\right),$$
(29)
$`F_\nu ^{(2)}`$ $``$ $`{\displaystyle _1^{\mathrm{}}}𝑑\gamma _2{\displaystyle _1^{\mathrm{}}}𝑑\gamma _1N_e(\gamma _2)N_e(\gamma _1)`$ (30)
$`\times R({\displaystyle \frac{3ϵ}{4\gamma _2^{2}}},\gamma _2)R({\displaystyle \frac{3}{4\gamma _2^{2}}}{\displaystyle \frac{3ϵ}{4\gamma _1^{2}}},\gamma _1)`$
$`\times F_{\mathrm{in}}\left({\displaystyle \frac{3}{4\gamma _2^{2}}}{\displaystyle \frac{3ϵ}{4\gamma _1^{2}}}\right)`$
with the definition $`F_{\nu ,\mathrm{in}}=(chϵ/2)n_{\mathrm{in}}(ϵ)F_{\mathrm{in}}(ϵ)`$. As the incident flux $`F_{\nu ,\mathrm{in}}`$ in the above expressions, the result from equation (10) should be used. The effects of low energy tail ($`\beta <1`$) in the electron distribution is neglected in performing the $`\gamma `$-integrals.
## 4. APPLICATION TO SAGITTARIUS A
In order to compare the resistive ADAF model with the current models of viscous ADAFs in their predictions of spectra from accretion flows, we apply the former model to Sgr A. The observed spectral data available so far have been compiled by Narayan et al. (1998). They assume that the interstellar column density is $`N_\mathrm{H}=6\times 10^{22}\mathrm{cm}^2`$ and the distance to the Galactic center is $`d=8.5\mathrm{kpc}`$. In judging the accuracy of fittings between the calculated and observed spectra, a considerable weight has been put on the high resolution data points such as the VLBI radio (86 GHz) data and the $`ROSAT`$ X-ray data. It should be kept in mind, however, that the $`ROSAT`$ data may be interpreted as an upper limit because its resolution (PSPC) $`20^{\prime \prime }`$ is not considered as satisfactory and that other issues like the value of $`N_\mathrm{H}`$ are still under discussion.
We discuss the two cases in the resistive model, which will be called the compact-disk and the extended-disk models, respectively. These names come from the difference in extension of the disk which is represented by the radius ratio of the inner to the outer edges, $`R_{\mathrm{in}}/R_{\mathrm{out}}=r_{\mathrm{in}}`$. As confirmed below, this value is largely affected by the choice of position of the inner edge, $`R_{\mathrm{in}}`$.
### 4.1. Compact Disk Model
According to the spirit of original resistive ADAF model, the inner edge in this case is determined by the magnetic flux conservation (K00). This gives an expression
$`R_{\mathrm{in}}=(1+\mathrm{\Delta }^1)^2R_{\mathrm{out}}\mathrm{\Delta }^2R_{\mathrm{out}},`$ (31)
where the last expression is valid only for thin disks ($`\mathrm{\Delta }1`$). Note that this procedure is independent of the notion of the marginally stable orbit around black holes. The outer edge has been fixed, on the other hand, from the mass conservation as
$`R_{\mathrm{out}}=\left({\displaystyle \frac{3GM\dot{M}^2}{B_0^4}}\right)^{1/5}.`$ (32)
Fig. 2 shows the best fit spectrum in this model and the set of best fit parameters is
$`M=3.9\times 10^5M_{},`$
$`\dot{M}=1.2\times 10^4\dot{M}_\mathrm{E}=1.0\times 10^6M_{}\mathrm{yr}^1,`$
$`|B_0|=0.7\mathrm{G},\mathrm{\Delta }=0.14\mathrm{rad}.`$ (33)
From these values, other quantities of our interest are fixed as follows:
$`R_{\mathrm{in}}=6.1\times 10^{12}\mathrm{cm},R_{\mathrm{out}}=3.1\times 10^{14}\mathrm{cm},`$
$`T=3.4\times 10^8r^1\mathrm{K},|b_\phi |=5.0r^1\mathrm{G},`$
$`\rho =1.8\times 10^{17}r^1\mathrm{g}\mathrm{cm}^3,`$
$`\tau _{es}=3.1\times 10^4.`$ (34)
Thus, it turns out that the inner edge of the present model is fairly large compared with the marginally stable orbit, $`R_{\mathrm{ms}}=3.5\times 10^{11}`$ cm, for a Schwarzschild hole of the above mass.
The changes in the spectrum caused by varying central mass $`M`$, accretion rate $`\dot{M}`$, external magnetic field $`|B_0|`$ and disk’s half-opening angle $`\mathrm{\Delta }`$ are demonstrated in Figs. 3, 4, 5 and 6, respectively. The spectral features are anyway quite analogous to those predicted by the viscous ADAF models. The results of a detailed comparison between the resistive and viscous ADAF models will be discussed in the final subsection, based on the predicted spectral features.
### 4.2. Extended Disk Model
In this model, the inner edge of the accretion disk is set at the radius of the marginally stable circular orbit around a Schwarzschild black hole,
$`R_{\mathrm{in}}=R_{\mathrm{ms}}=3R_\mathrm{G}={\displaystyle \frac{6GM}{c^2}},`$ (35)
where $`R_\mathrm{G}`$ is the gravitational radius of the hole. This choice is motivated by the expectation that at around this radius the infall velocity inevitably becomes of the order of the rotational velocity, (i.e., $`\mathrm{}1`$ where $`\mathrm{}`$ is the magnetic Reynolds number, see K99, K00). The definition of the outer edge is the same as in the compact disk model. Note that the above definition of the inner edge is adopted also in the viscous ADAF models.
The best fit parameters in this model are
$`M=1.0\times 10^6M_{},`$
$`\dot{M}=1.3\times 10^4\dot{M}_\mathrm{E}=2.9\times 10^6M_{}\mathrm{yr}^1,`$
$`|B_0|=1.0\times 10^6\mathrm{G},\mathrm{\Delta }=0.20\mathrm{rad}.`$ (36)
These are used to fix the values of various scaled quantities:
$`R_{\mathrm{in}}=8.9\times 10^{11}\mathrm{cm},R_{\mathrm{out}}=2.7\times 10^{19}\mathrm{cm},`$
$`T=1.0\times 10^4r^1\mathrm{K},|b_\phi |=5.0\times 10^6r^1\mathrm{G},`$
$`\rho =5.9\times 10^{25}r^1\mathrm{g}\mathrm{cm}^3,`$
$`\tau _{es}=1.3\times 10^6.`$ (37)
The best fit curve is shown in Fig. 7. The changes in the spectrum caused by varying central mass $`M`$, accretion rate $`\dot{M}`$, external magnetic field $`|B_0|`$ and disk’s half-opening angle $`\mathrm{\Delta }`$ are demonstrated in Figs. 8, 9, 10 and 11, respectively. The spectral shapes are very different from those of the compact-disk case and of the viscous ADAF models. Synchrotron emission has a very wide peak and bremsstrahlung is negligibly small. The former fact is due to a high temperature at the inner edge (see sub-subsection 4.3.2) and the latter, to lower densities in the disk. The emission in the X-ray band is supported by the inverse Compton scattering from the radio band. The temperature near the outer edge falls even to such a small value that the assumption of complete ionization becomes invalid. Although the position of outer edge may seem to be irrelevant from a viewpoint of spectrum, it is nevertheless important also in this case as a fitting boundary of the inner magnetic field to the external one. The fitting predicts that the boundary value is comparable to the interstellar field (a few $`\mu `$G).
The fitting both to 86 GHz and ROSAT X-ray data points is possible also in this model. However, it is clear that the fitting curve runs above the observed upper limits in the IR band. The fitting in the frequency range from 100 to 1000 GHz also becomes considerably poor compared with the case of compact disk. For these reasons, we judge that this model cannot reproduce the observed broadband spectrum of Sgr A. This fact suggests again that the inner edge of the accretion disk does not coincide with the marginally stable orbit. The wide range of the disk’s radii which is obtained from this fitting implies that $`\mathrm{}(R_{\mathrm{out}})6\times 10^3`$. Since $`\mathrm{}(R)`$ represents the ratio of toroidal to poloidal magnetic fields, most parts of the disk are very likely to be unstable to global MHD instabilities of helical type. For this reason too, we consider that the present case (i.e., $`R_{\mathrm{in}}=R_{\mathrm{ms}}`$) is quite unrealistic, at least, for Sgr A.
### 4.3. Viscous v.s. Resistive ADAFs
#### 4.3.1 Dependence on Black Hole Mass
The spectra calculated from ADAF models of both viscous and resistive types commonly have the saturated part at the lower ends of the spectra due to the synchrotron self-absorption. It is of great interest to see that the luminosity $`\nu L_\nu `$ of this part is essential to determine the mass of the central black hole, in both types of models. Especially, in the viscous model, the luminosity of this frequency part is determined almost only by the black hole mass. The reason is as follows.
The temperature in ADAFs may be considered essentially as the ion virial temperature and hence decreases as $`R^1`$. Apart from a numerical factor due to a reduced Keplerian rotation, this is exactly true in the resistive model. This is also true in the viscous models for the main part of an accretion flow except in the inner region where the electron temperature deviates from the ion temperature and remains almost constant (e.g., Narayan & Yi 1995b). Therefore, the contribution to the spectrum from each annulus of radius $`R`$ and width d$`R`$ is equal. Integrating these contributions up to the outer edge, we obtain $`L_\nu ^{\mathrm{RJ}}T_e(R_{\mathrm{in}})R_{\mathrm{in}}R_{\mathrm{out}}=T_e(R_{\mathrm{out}})R_{\mathrm{out}}^2`$, where $`R_{\mathrm{in}}`$ is the radius of the disk’s inner edge in the resistive model and of the outer edge of the two-temperature region in the viscous models.
We have $`T_eRm`$ commonly to both types of ADAF models. Further, since radius scales as the gravitational radius in the case of viscous ADAFs, we obtain the mass dependence
$$L_\nu ^{\mathrm{RJ}}m^2\text{(viscous ADAF)},$$
(38)
confirming the above statement. On the other hand, in the case of resistive ADAFs, we have
$`L_\nu ^{\mathrm{RJ}}b_0^{4/5}\dot{m}^{2/5}m^{8/5}\text{(resistive ADAF)},`$ (39)
where the dependences on the parameters other than $`m`$ have come from the expression of $`R_{\mathrm{out}}`$. In spite of these dependences, the mass dependence is essential also in this case. This is because the dependence on $`\dot{m}`$ is rather weak and the value of $`b_0`$ is strongly restricted from the position of the synchrotron peak (see the discussion below).
#### 4.3.2 Synchrotron Peak
We estimate the synchrotron peak frequency following Mahadevan (1997), and examine its behavior in both viscous and resistive models. For each annulus of radius $`R`$ and width d$`R`$, the synchrotron photons in the radio range up to a critical frequency $`\nu _\mathrm{c}`$ are strongly self-absorbed and result in the Rayleigh-Jeans spectrum. Therefore, the critical frequency of the spectrum is determined by equating the contributions to $`L_\nu `$ from optically thick and thin sides of the frequency:
$`2\pi {\displaystyle \frac{\nu _\mathrm{c}^2}{c^2}}k_\mathrm{B}T_e(R)2\pi R\mathrm{d}R`$
$`=4.43\times 10^{30}{\displaystyle \frac{4\pi n_e\nu _\mathrm{c}}{K_2(1/\theta _e)}}I^{}(x_\mathrm{c})4\pi \mathrm{\Delta }R^2\mathrm{d}R,`$ (40)
where $`x_\mathrm{c}`$ is defined as $`x_\mathrm{c}=2\nu _\mathrm{c}/(3\nu _0\theta _e^2)`$. Solving this equation, we can determine the value of $`x_\mathrm{c}`$ numerically (Appendix B of Mahadevan 1997). Provided that this value does not depend strongly on $`R`$, $`\mathrm{\Delta }`$ and other parameters, we obtain
$`\nu _\mathrm{c}={\displaystyle \frac{3}{2}}\theta _e^2\nu _0x_\mathrm{c}T_e^{2}(r)B(r).`$ (41)
If the disk has uniform temperature and magnetic field, then the synchrotron peak is rather sharp and has a well-defined peak frequency at $`\nu _\mathrm{c}`$.
When they vary with the radius $`R`$, however, substitution of the $`r`$-dependences of $`T_\mathrm{e}`$ and $`B`$ in both viscous and resistive ADAF models yield
$`\nu _\mathrm{c}`$ $``$ $`\alpha ^{1/2}(1\beta )^{1/2}\dot{m}^{1/2}m^{1/2}r^{13/4}`$ (42)
$`\text{(viscous ADAF)},`$
$``$ $`\delta ^1b_0^{13/5}\dot{m}^{4/5}m^{4/5}r^3`$
$`\text{(resistive ADAF)}.`$
This means that $`\nu _\mathrm{c}`$ is larger for smaller radii and the higher most cutoff is due to the inner edge. The position of peak of the superposed emission is then given as $`\nu _\mathrm{p}=\nu _\mathrm{c}(r_\mathrm{p})`$, where $`r_\mathrm{p}`$ is the radius whose contribution to the synchrotron emission is most dominant. The fairly narrow peak obtained in the compact-disk case indicates that $`r_\mathrm{p}`$ is located near the inner edge and the global peak shape is determined mainly by the inner most region of the disk.
On the contrary, the synchrotron peak becomes very broad and dull in the extended-disk case. We have confirmed that the low-frequency side of the broad peak is due to a superposition of the contributions from annuli of $`R_{\mathrm{ms}}10R_{\mathrm{ms}}`$. However, the dull shape on the high-frequency side of the peak may be mainly due to a resulting high temperature ($`T3\times 10^{10}`$ K) at the smaller inner edge. Actually, owing to this high temperature and low densities near the inner edge, the synchrotron self-absorption becomes less important in the high-frequency radio band and the intrinsic shape of the synchrotron emission at the mildly relativistic temperature (Mahadevan et al. 1996) can appear on the high-frequency side.
In any case, since $`r_\mathrm{p}`$ is a numerical factor, we can speak of the parameter dependences of the peak frequency $`\nu _\mathrm{p}`$ based on equation (42). Note that the dependences on $`m`$ and $`\dot{m}`$ have different senses in the different ADAFs. The most important difference between the two models is that the dependence of $`\nu _\mathrm{p}`$ on the magnetic field is much more sensitive in the resistive model. Therefore, the field strength is determined more accurately there. All the predicted dependences on $`m`$, $`\dot{m}`$, $`b_0`$ and $`\delta `$ are qualitatively confirmed in Figs. 2 through 5. From the above considerations on the synchrotron peak, we think that the improvement of observational quality in submillimeter range is most important for obtaining more exact values of the disk parameters.
#### 4.3.3 Bremsstrahlung
We shall try here to grasp the qualitative behavior of the contribution from bremsstrahlung according to the usual non-relativistic scheme. The contribution to a given frequency $`\nu `$ from optically thin plasma in an annular volume of width d$`R`$ is proportional to $`\rho ^2T^{1/2}\mathrm{exp}[h\nu /k_\mathrm{B}T]R^2\mathrm{d}R`$. Apart from the exponential factor, we have
$`\rho ^2T^{1/2}R^2\mathrm{d}R`$ $``$ $`\alpha ^2\dot{m}^2mr^{1/2}\mathrm{d}r\text{(viscous ADAF)},`$ (43)
$``$ $`\delta ^4b_0^{2/5}\dot{m}^{11/5}m^{4/5}r^{1/2}\mathrm{d}r`$
$`\text{(resistive ADAF)}.`$
Therefore, the relative importance of the inner and outer parts of a disk can be seen from the ratio,
$$f\frac{\rho _{\mathrm{in}}^2T_{\mathrm{in}}^{1/2}R_{\mathrm{in}}^2\mathrm{exp}[h\nu /k_\mathrm{B}T_{\mathrm{in}}]}{\rho _{\mathrm{out}}^2T_{\mathrm{out}}^{1/2}R_{\mathrm{out}}^2\mathrm{exp}[h\nu /k_\mathrm{B}T_{\mathrm{out}}]}\zeta ^{\pm 1/2}\mathrm{exp}[h\nu /k_\mathrm{B}T_{\mathrm{out}}],$$
(44)
where $`\zeta R_{\mathrm{out}}/R_{\mathrm{in}}=r_{\mathrm{in}}^1`$, and the upper and lower signs in its exponent are for the viscous and resistive ADAFs, respectively.
It is evident from the above ratio that, in the viscous ADAF, the contribution from the inner disk is always dominant (i.e., $`f>1`$) irrespective of the frequency $`\nu `$. In the resistive ADAF, however, it depends on the frequency, so that we introduce the critical frequency $`\nu _\mathrm{c}^{}`$ by the relation $`f=1`$. This yields
$$\nu _\mathrm{c}^{}=\frac{\mathrm{ln}\zeta }{2}\frac{k_\mathrm{B}T_{\mathrm{out}}}{h}.$$
(45)
Then, the inner part contributes to the frequency range $`\nu >\nu _\mathrm{c}^{}`$ and the outer part, to $`\nu <\nu _\mathrm{c}^{}`$. In fact, the critical frequency roughly coincides with the peak frequency of the bremsstrahlung. The luminosity above $`\nu _\mathrm{c}^{}`$ can be roughly estimated, by putting $`\mathrm{d}rrr_{\mathrm{in}}`$, as
$`L_\nu ^{\mathrm{br}}`$ $``$ $`\alpha ^2\dot{m}^2m\text{(viscous ADAF)},`$ (46)
$``$ $`\delta ^1b_0^{2/5}\dot{m}^{11/5}m^{4/5}\text{(resistive ADAF)},`$
because $`r_{\mathrm{in}}`$ is a numerical constant and, in particular, equal to $`\delta ^2`$ in the resistive ADAFs.
In the viscous ADAF models and in the compact-disk case of the resistive ADAF, the contributions from bremsstrahlung cause an X-ray bump in each predicted spectrum. The dependences of this peak on the parameters $`m`$ and $`\dot{m}`$ are qualitatively confirmed in Figs. 2 and 3, but those on $`b_0`$ and $`\delta `$ are somewhat different from the above prediction, indicating a limitation of such a crude estimate as the above. The critical frequencies calculated from the best fit values for the compact and extended disks are $`1.4\times 10^{19}`$ Hz and $`1.8\times 10^{15}`$ Hz, respectively. The former value is in good agreement with the peak of the reproduced spectrum. In the extended-disk case, the contribution from the bremsstrahlung is negligibly small because the density throughout the disk becomes too small, and the X-ray range of the spectrum is explained by the once-scattered Compton photons.
## 5. Summary $`\&`$ Discussion
To summarize the examinations in the previous section, both viscous and resistive ADAF models can explain the observed spectrum of Sgr A equally well. In spite of large differences in the basic mechanisms working in both models, the calculated spectra are quite similar, except for the extended-disk case in the resistive model. This fact suggests that also the resistive ADAF model is quite powerful in explaining the behavior of other low luminosity AGNs (Narayan, Mahadevan & Quataert 1998). In addition to these analogous aspects, the resistive model seems to have a possibility to explain such an essentially different situation as appeared in the extended disk case. In any case, when the presence of an ordered magnetic field should be taken seriously in some AGNs or in some stellar-size black holes then the resistive ADAF model, whose predictions on the radiation spectra are examined in this paper, will serve the purpose.
One of the most remarkable features of the ADAF models is that the mass of the central black hole seems to be determined only from the fitting to the self-absorbed part of the observed spectrum. In the case of Sgr A, the resistive ADAF model (hereafter restricting to the case of compact disk) predicts the central mass of $`3.9\times 10^5M_{}`$ while the viscous models predict $`1.0\times 10^6M_{}`$ (Manmoto et.al 1997) and $`2.5\times 10^6M_{}`$ (Narayan et al. 1998). The accuracy of the fittings for other disk parameters than the black hole mass will be greatly improved by the precise determination of the position and height of the synchrotron peak from observations.
The black hole mass predicted by the resistive ADAF model is evidently smaller compared with the predictions of the viscous ADAF models. The latter values are consistent with the dynamically reduced value of $`2.5\times 10^6M_{}`$ (Haller et al. 1996; Eckart & Genzel 1997), which may be considered as an upper limit for the black hole mass. In the history of viscous ADAF models, the predicted black-hole mass was as small as $`7\times 10^5M_{}`$ (Narayan et al. 1995). Afterwards, by the inclusion of compressive heating, it becomes consistent with the dynamical mass. Since this change is mainly due to the decrease in electron temperature (Narayan et al. 1998), the prediction of the resistive ADAF model may also be increased if the development of its two-temperature versions results in a lower electron temperature. As for the compressive heating, it is already included in the resistive model.
In spite of the resemblance in the predicted spectral shape, there are of course many differences in the predictions of the viscous and resistive models. The precise dependences on the relevant quantities of the luminosities of the self-absorbed part, the synchrotron peak and the X-ray bump are different. Especially, the dependence of the synchrotron peak-frequency on the strength of magnetic field is much stronger for the resistive model. The essential difference in the geometry of an accretion flow may be in the radius of the inner edge rather than in its vertical thickness. The prediction of the resistive ADAF model for the inner-edge radius of the disk around Sgr A is $`20R_{\mathrm{ms}}`$, instead of the radius of marginally stable circular orbit $`R_{\mathrm{ms}}`$. Although this result justifies the neglection of the general relativistic effects in our treatment, various questions may be raised about the behavior of infalling plasmas. As for this point we only present an idea below from a viewpoint of global consistency, because its detailed analyses are beyond the scope of this paper.
Fig. 1 shows an overview of the flow and magnetic field configurations (see K00, for more details). The accretion flow would be decelerated near the inner edge by the presence of a strong poloidal magnetic field which is maintained by the sweeping effect of the flow. As a result, a certain fraction of the accreting plasma will be turned its direction to go along the poloidal field lines, although the remaining fraction may fall into the central black hole. If the poloidal current driven in the accretion disk can close its circuit successfully around distant regions and along the polar axis, a set of bipolar jets will be formed (Kaburaki $`\&`$ Itoh 1987). Even if the mechanism for formation of jet does not work well, the plasma within the inner edge is likely to extend to the polar regions.
The presence of the plasma within the inner edge of an accretion disk and near the polar axis can be a possible source of the excess above the self-absorbed slope in radio band of the observed spectrum. Very recent VLBI observations of Sgr A (Krichbaum et al. 1998; Lo et al. 1998) report that its intrinsic sizes in the east-west direction at 215 GHz and 68GHz are about 20 $`R_\mathrm{G}`$ (with $`M=2.5\times 10^6M_{}`$). A half of this size (i.e., its radius) is just comparable to the size of the inner edge $`60R_\mathrm{G}`$ of our model fitting with $`M=4\times 10^5M_{}`$. However, the value of the black hole mass estimated from the spectral fitting may be increased if there is a possibility for Sgr A to have a wind-type mass loss from the surfaces of its disk (such possibilities have been noted for various types of objects by, e.g., Blandford & Begelman 1999; Di Matteo et al. 1999; Quataert & Narayan 1999). In such a case, the VLBI component becomes smaller than the size of the inner edge.
From the standpoint of the resistive ADAF, therefore, the above observations should be interpreted as suggesting the presence of a compact structure which is comparable to or smaller than the inner-edge radius of the accretion disk. In view of the vertical elongation of this component reported by Lo et al. (1998), this structure is very likely to be the root of a jet as suggested by them. This picture is very consistent with the view described above in relation to Fig. 1. In this case, however, the location of the self-absorbed slope in the $`\nu `$-$`\nu L_\nu `$ diagram should be slightly shifted towards the higher-frequency side so that the VLBI data points can be regarded as an excess from the disk’s contribution.
One of the authors (M. K.) would like to thank Tadahiro Manmoto for many valuable comments on the viscous ADAF models. He is also grateful to Umin Lee for his suggestions on some numerical technics. |
warning/0002/hep-ph0002101.html | ar5iv | text | # 𝑍-peak subtracted representation of Bhabha scattering and search for new physics effects
## I Introduction.
A convenient way of searching for virtual new physics effects in $`e^+e^{}`$ annihilation processes has been recently proposed; it consists of the use of the so-called ”Z-peak subtracted representation” that allows to take automatically into account the severe constraints imposed by the high precision measurements performed at the $`Z`$ peak by LEP1 and SLC. This is achieved by choosing as ”theoretical inputs” the measured values of the partial $`Z`$ widths $`\mathrm{\Gamma }_f`$ and of the effective weak angle $`sin^2\theta _{\mathrm{eff}}`$ together with $`\alpha (0)`$. For each $`e^+e^{}f\overline{f}`$ ($`fe`$) annihilation process, all the one-loop standard model (SM) or new physics (NP) effects are described by four functions of the energy ($`\sqrt{q^2}`$) and of the scattering angle $`\theta `$ which are subtracted at $`q^2=M_Z^2`$ or at $`q^2=0`$ in order to take the inputs into account. We defer to ref. for further details.
This procedure is especially suitable for the search of NP effects which grow with the energy. Various applications were made for Supersymmetry , Technicolour , anomalous gauge couplings, higher $`Z^{}`$ bosons, 4-fermion contact terms and extra-dimensions . The description was shown to be particularly useful when NP is characterized by an effective scale $`\mathrm{\Lambda }`$ which is much higher than the actual energy range of $`\sqrt{q^2}`$. In the particular case of universal $`\theta `$ independent effects all the information on NP can be then conveyed into three constants called, $`\delta _{Z,s,\gamma }`$, that can be viewed as the generalization, beyond the $`Z`$ peak, of the $`T`$, $`S`$ or $`ϵ_{1,3}`$ description.
This representation was not yet applied to Bhabha scattering because of the complication generated by the presence of $`t`$-channel photon and $`Z`$ exchanges. The purpose of the present paper is to fill this lack. We shall show here in Section 2 that the whole $`Z`$-peak subtracted formalism can be extended to Bhabha scattering in a very natural way. This will allow to use this process, together with the other $`e^+e^{}f\overline{f}`$ ones, in order to improve the constraints on possible NP contributions, as illustrated in Section 3. In Section 4, we shall show in the numerical applications that the gain thus obtained can be substancial.
## II The $`Z`$-peak subtracted representation
We shall first summarize the results of the $`Z`$-peak subtracted representation, described in previous dedicated papers , by writing the general $`e^+e^{}f\overline{f}`$ ($`fe`$) scattering amplitude at one loop as the sum of an effective photon and an effective $`Z`$ amplitude with couplings $`g_{Vj}^\gamma (q^2,\theta )`$, $`g_{Vj}^Z(q^2,\theta )`$, $`g_{Aj}^Z(q^2,\theta )`$, where the index $`j`$ denotes either the initial electron ($`j=e`$) or the final fermion ($`j=fe`$).
$`𝒜(q^2,\theta )={\displaystyle \frac{i}{q^2}}\overline{v}(e^+)\gamma ^\mu g_{Ve}^{(\gamma )}(q^2,\theta )u(e^{}).\overline{u}(f)\gamma _\mu g_{Vf}^{(\gamma )}(q^2,\theta )v(\overline{f})+{\displaystyle \frac{i}{q^2M_Z^2+iM_Z\mathrm{\Gamma }_Z}}.`$ (1)
$`\overline{v}(e^+)\gamma ^\mu [g_{Ve}^{(Z)}(q^2,\theta )g_{Ae}^{(Z)}(q^2,\theta )\gamma ^5]u(e^{}).\overline{u}(f)\gamma _\mu [g_{Vf}^{(Z)}(q^2,\theta )g_{Af}^{(Z)}(q^2,\theta )\gamma ^5]v(\overline{f})`$ (2)
The aforementioned inputs are taken into account by imposing that the total amplitude takes the required value at $`q^2=0`$ and at $`q^2=M_Z^2`$. Following the method of , this amounts to use a subtraction procedure which allows to write:
$`g_{Ve}^\gamma (q^2,\theta )=\sqrt{4\pi \alpha (0)}Q_e[1+{\displaystyle \frac{1}{2}}\stackrel{~}{\mathrm{\Delta }}_{\alpha ,ef}(q^2,\theta )]`$ (3)
$`g_{Vf}^\gamma (q^2,\theta )=\sqrt{4\pi \alpha (0)}Q_f[1+{\displaystyle \frac{1}{2}}\stackrel{~}{\mathrm{\Delta }}_{\alpha ,ef}(q^2,\theta )]`$ (4)
$`g_{Ae}^\gamma (q^2,\theta )=g_{Af}^\gamma (q^2,\theta )=0`$ (5)
$`g_{Ve}^Z=\gamma _e^{\frac{1}{2}}I_{3e}\stackrel{~}{v}_e[1{\displaystyle \frac{1}{2}}R_{ef}(q^2,\theta ){\displaystyle \frac{4\stackrel{~}{s}_e\stackrel{~}{c}_e}{\stackrel{~}{v}_e}}|Q_f|V_{ef}^{\gamma Z}(q^2,\theta )]`$ (6)
$`g_{Vf}^Z(q^2,\theta )=\gamma _f^{\frac{1}{2}}I_{3f}\stackrel{~}{v}_f[1{\displaystyle \frac{1}{2}}R_{ef}(q^2,\theta ){\displaystyle \frac{4\stackrel{~}{s}_e\stackrel{~}{c}_e}{\stackrel{~}{v}_f}}|Q_f|V_{ef}^{Z\gamma }(q^2,\theta )]`$ (7)
$`g_{Ae}^Z(q^2,\theta )=\gamma _e^{\frac{1}{2}}I_{3e}[1{\displaystyle \frac{1}{2}}R_{ef}(q^2,\theta )]`$ (8)
$`g_{Af}^Z(q^2,\theta )=\gamma _f^{\frac{1}{2}}I_{3f}[1{\displaystyle \frac{1}{2}}R_{ef}(q^2,\theta )]`$ (9)
(10)
with the $`Z`$-peak inputs
$$\gamma _j^{\frac{1}{2}}=[\frac{48\pi \mathrm{\Gamma }_j}{N_jM_Z(1+\stackrel{~}{v}_j^2)}]^{\frac{1}{2}}$$
(11)
and
$$\stackrel{~}{v}_j=14|Q_j|\stackrel{~}{s}_j^2$$
(12)
where $`\stackrel{~}{s}_j^2=1\stackrel{~}{c}_j^2`$ is the weak effective angle measured through the forward-backward or polarization asymmetries in the final channel $`j`$, $`\stackrel{~}{s}_e\stackrel{~}{s}_\mu \stackrel{~}{s}_\tau `$ and $`N_j`$ is the colour factor with QCD corrections at $`Z`$-peak.
The quantities $`\stackrel{~}{\mathrm{\Delta }}_{\alpha ,ef}(q^2,\theta )`$, $`R_{ef}(q^2,\theta )`$, $`V_{ef}^{\gamma Z}(q^2,\theta )`$, $`V_{ef}^{Z\gamma }(q^2,\theta )`$ contain all the $`q^2,\theta `$ dependent parts of the scattering amplitude due to SM or NP at one-loop. In our approach, they consist of certain finite combinations of self-energies, vertices and boxes that are automatically gauge independent.
For an additional four fermion amplitude
$$\overline{v}(e^+)\gamma ^\mu [a(q^2,\theta )b(q^2,\theta )\gamma ^5]u(e^{}).\overline{u}(f)\gamma _\mu [c(q^2,\theta )d(q^2,\theta )\gamma ^5]v(f)$$
(13)
where $`a`$, $`b`$, $`c`$ and $`d`$ are $`𝒪(\alpha )`$ quantities, one easily obtains the corresponding projections on the photon and $`Z`$ Lorentz structures:
$`\stackrel{~}{\mathrm{\Delta }}_{\alpha ,ef}(q^2,\theta )=q^2{\displaystyle \frac{[a(q^2,\theta )b(q^2,\theta )\stackrel{~}{v}_e][c(q^2,\theta )d(q^2,\theta )\stackrel{~}{v}_f]}{e^2Q_eQ_f}}`$ (14)
$`R_{ef}(q^2,\theta )=(q^2M_Z^2){\displaystyle \frac{4\stackrel{~}{s}_e^2\stackrel{~}{c}_e^2b(q^2,\theta )d(q^2,\theta )}{e^2I_{3e}I_{3f}}}`$ (15)
$`V_{ef}^{\gamma Z}(q^2,\theta )=(q^2M_Z^2){\displaystyle \frac{[a(q^2,\theta )b(q^2,\theta )\stackrel{~}{v}_e]2\stackrel{~}{s}_e\stackrel{~}{c}_ed(q^2,\theta )}{e^2Q_eI_{3f}}}`$ (16)
$`V_{ef}^{Z\gamma }(q^2,\theta )=(q^2M_Z^2){\displaystyle \frac{[c(q^2,\theta )d(q^2,\theta )\stackrel{~}{v}_f]2\stackrel{~}{s}_e\stackrel{~}{c}_eb(q^2,\theta )}{e^2Q_fI_{3e}}}`$ (17)
We have given in ref. and in the Appendix of ref. the expression of the general polarized $`e^+e^{}f\overline{f}`$ differential cross section in terms of these four functions. From this one obtains, for example, the integrated cross sections $`\sigma _f`$ and the asymmetries $`A_{FB,f}`$, $`A_{LR,f}`$.
To generalize our approach to the case $`f=e`$, Bhabha scattering, it is convenient to write the scattering amplitude at one loop as the sum of two (s-channel and t-channel) components:
$$𝒜_{ee}=𝒜_s(q^2,\theta )+𝒜_t(q^2,\theta )$$
(18)
The procedure that we have illustrated in the $`e^+e^{}\overline{f}f`$ ($`fe`$) case applies directly to the $`s`$-channel part of the Bhabha amplitude. In this case we can drop the index $`ef`$, as $`ef`$ and we have $`V^{\gamma Z}(q^2,\theta )V^{Z\gamma }(q^2,\theta )V(q^2,\theta )`$, so that we only deal with three independent functions $`\stackrel{~}{\mathrm{\Delta }}_\alpha (q^2,\theta )`$, $`R(q^2,\theta )`$ and $`V(q^2,\theta )`$.
It is now straightforward to check that the same procedure can be applied step by step to the $`t`$-channel component. In full generality the latter can be written as
$`𝒜_t(q^2,\theta )={\displaystyle \frac{i}{t}}\overline{v}(e^+)\gamma ^\mu \overline{g}_{Ve}^{(\gamma )}(q^2,\theta )v(e^+).\overline{u}(e^{})\gamma _\mu \overline{g}_{Vf}^{(\gamma )}(q^2,\theta )u(e^{})+{\displaystyle \frac{i}{tM_Z^2}}.`$ (19)
$`\overline{v}(e^+)\gamma ^\mu [\overline{g}_{Ve}^{(Z)}(q^2,\theta )\overline{g}_{Ae}^{(Z)}(q^2,\theta )\gamma ^5]v(e^+)\overline{u}(e^{})\gamma _\mu [\overline{g}_{Vf}^{(Z)}(q^2,\theta )\overline{g}_{Af}^{(Z)}(q^2,\theta )\gamma ^5]u(e^{})`$ (20)
(21)
with the $`t`$-channel effective couplings:
$`\overline{g}_{Ve}^\gamma (q^2,\theta )=\sqrt{4\pi \alpha (0)}Q_e[1+{\displaystyle \frac{1}{2}}\overline{\stackrel{~}{\mathrm{\Delta }}}_\alpha (q^2,\theta )]`$ (22)
$`\overline{g}_{Ve}^Z(q^2,\theta )=\gamma _e^{\frac{1}{2}}I_{3e}\stackrel{~}{v}_e[1{\displaystyle \frac{1}{2}}\overline{R}(q^2,\theta ){\displaystyle \frac{4\stackrel{~}{s}_e\stackrel{~}{c}_e}{\stackrel{~}{v}_e}}|Q_f|\overline{V}(q^2,\theta )]`$ (23)
$`\overline{g}_{Ae}^Z(q^2,\theta )=\gamma _e^{\frac{1}{2}}I_{3e}[1{\displaystyle \frac{1}{2}}\overline{R}(q^2,\theta )]`$ (24)
(25)
where the new functions $`\overline{\stackrel{~}{\mathrm{\Delta }}}_\alpha (q^2,\theta )`$, $`\overline{R}(q^2,\theta )`$ and $`\overline{V}(q^2,\theta )`$ are obtained from the previous $`s`$-channel ones (without bar) by ($`q^2t`$) crossing relations:
$$\overline{\stackrel{~}{\mathrm{\Delta }}}_\alpha (q^2,\theta )=[\stackrel{~}{\mathrm{\Delta }}_\alpha (q^2,\theta )](q^2t=\frac{q^2}{2}(1cos\theta );cos\theta 1+\frac{2q^2}{t})$$
(26)
and analogously for $`\overline{R}`$, $`\overline{V}`$.
The general expression of the polarized Bhabha differential cross section obtained from the sum of the $`s`$-channel (2) and $`t`$-channel (21) amplitudes is given in the Appendix, in the form:
$$\frac{d\sigma }{dcos\theta }=(1PP^{})\frac{d\sigma ^1}{dcos\theta }+(1+PP^{})\frac{d\sigma ^2}{dcos\theta }+(P^{}P)\frac{d\sigma ^P}{dcos\theta }$$
(27)
where $`P`$ and $`P^{}`$ are the initial $`e^{}`$, $`e^+`$ polarizations.
We shall write the three differential cross sections as the sum of a Born term and a one loop contribution, $`d\sigma ^i=(d\sigma ^i)^{Born}+(d\sigma ^i)^{(1)}`$.
The Born term is given by
$`({\displaystyle \frac{d\sigma ^1}{dcos\theta }})^{Born}={\displaystyle \frac{\pi \alpha ^2}{q^2}}\{{\displaystyle \frac{t^2+u^2}{q^4}}+{\displaystyle \frac{u^2}{t^2}}+{\displaystyle \frac{2u^2}{tq^2}}`$ (28)
$`+2({\displaystyle \frac{3\mathrm{\Gamma }_e}{\alpha M_Z}})[{\displaystyle \frac{(u^2t^2+\stackrel{~}{v}_e^2(u^2+t^2))(q^2M_Z^2)}{(1+\stackrel{~}{v}_e^2)q^2[(q^2M_Z^2)^2+M_Z^2\mathrm{\Gamma }_Z^2]}}+{\displaystyle \frac{u^2}{q^2(tM_Z^2)}}+{\displaystyle \frac{u^2(q^2M_Z^2)}{t[(q^2M_Z^2)^2+M_Z^2\mathrm{\Gamma }_Z^2]}}`$ (29)
$`+{\displaystyle \frac{u^2}{t(tM_Z^2)}}]+({\displaystyle \frac{3\mathrm{\Gamma }_e}{\alpha M_Z}})^2[{\displaystyle \frac{(t^2+u^2)(1+\stackrel{~}{v}_e^2)^2+4\stackrel{~}{v}_e^2(u^2t^2)}{(1+\stackrel{~}{v}_e^2)^2[(q^2M_Z^2)^2+M_Z^2\mathrm{\Gamma }_Z^2]}}`$ (30)
$`+{\displaystyle \frac{u^2[(1+\stackrel{~}{v}_e^2)^2+4\stackrel{~}{v}_e^2]}{(1+\stackrel{~}{v}_e^2)^2(tM_Z^2)}}({\displaystyle \frac{2(q^2M_Z^2)}{[(q^2M_Z^2)^2+M_Z^2\mathrm{\Gamma }_Z^2]}}+{\displaystyle \frac{1}{(tM_Z^2)}})]\}`$ (31)
$`({\displaystyle \frac{d\sigma ^2}{dcos\theta }})^{Born}=\pi \alpha ^2q^2\{{\displaystyle \frac{1}{t^2}}+2({\displaystyle \frac{3\mathrm{\Gamma }_e}{\alpha M_Z}}){\displaystyle \frac{(\stackrel{~}{v}_e^21)}{(1+\stackrel{~}{v}_e^2)t(tM_Z^2)}}`$ (32)
$`+({\displaystyle \frac{3\mathrm{\Gamma }_e}{\alpha M_Z}})^2{\displaystyle \frac{(1\stackrel{~}{v}_e^2)^2}{(1+\stackrel{~}{v}_e^2)^2(tM_Z^2)^2}}\}`$ (33)
$`({\displaystyle \frac{d\sigma ^P}{dcos\theta }})^{Born}={\displaystyle \frac{4\pi \alpha ^2u^2}{q^2}}[{\displaystyle \frac{\stackrel{~}{v}_e}{(1+\stackrel{~}{v}_e^2)}}]\{({\displaystyle \frac{3\mathrm{\Gamma }_e}{\alpha M_Z}})({\displaystyle \frac{1}{q^2}}+{\displaystyle \frac{1}{t}})[{\displaystyle \frac{(q^2M_Z^2)}{[(q^2M_Z^2)^2+M_Z^2\mathrm{\Gamma }_Z^2]}}+{\displaystyle \frac{1}{(tM_Z^2)}}]`$ (34)
$`+({\displaystyle \frac{3\mathrm{\Gamma }_e}{\alpha M_Z}})^2[{\displaystyle \frac{1}{[(q^2M_Z^2)^2+M_Z^2\mathrm{\Gamma }_Z^2]}}(1+{\displaystyle \frac{2(q^2M_Z^2)}{(tM_Z^2)}})+{\displaystyle \frac{1}{(tM_Z^2)^2}}]\}`$ (35)
and the one loop contributions of the three functions can be written in the condensed way:
$`({\displaystyle \frac{d\sigma ^1}{dcos\theta }})^{(1)}={\displaystyle \frac{\pi \alpha ^2}{q^2}}\{(t^2+u^2)G_1(q^2,q^2)+(u^2t^2)G_2(q^2,q^2)`$ (36)
$`+u^2[G_1(t,t)+G_2(t,t)+2G_1(q^2,t)+2G_2(q^2,t)]\}`$ (37)
$`({\displaystyle \frac{d\sigma ^2}{dcos\theta }})^{(1)}=\pi \alpha ^2q^2[G_1(t,t)G_2(t,t)]`$ (38)
$`({\displaystyle \frac{d\sigma ^P}{dcos\theta }})^{(1)}={\displaystyle \frac{4\pi \alpha ^2u^2}{q^2}}[G_3(q^2,q^2)+G_3(q^2,t)+G_3(t,q^2)+G_3(t,t)]`$ (39)
with
$`G_1(x,y)={\displaystyle \frac{\stackrel{~}{\mathrm{\Delta }}_\alpha (x)+\stackrel{~}{\mathrm{\Delta }}_\alpha (y)}{xy}}`$ (40)
$`+({\displaystyle \frac{3\mathrm{\Gamma }_e}{\alpha M_Z}}){\displaystyle \frac{\stackrel{~}{v}_e^2}{(1+\stackrel{~}{v}_e^2)}}[{\displaystyle \frac{\stackrel{~}{\mathrm{\Delta }}_\alpha (x)R(y)\frac{8\stackrel{~}{s}_e\stackrel{~}{c}_e}{\stackrel{~}{v}_e}V(y)}{x(yM_Z^2)}}+{\displaystyle \frac{\stackrel{~}{\mathrm{\Delta }}_\alpha (y)R(x)\frac{8\stackrel{~}{s}_e\stackrel{~}{c}_e}{\stackrel{~}{v}_e}V(x)}{y(xM_Z^2)}}]`$ (41)
$`({\displaystyle \frac{3\mathrm{\Gamma }_e}{\alpha M_Z}})^2[{\displaystyle \frac{R(x)+R(y)+\frac{8\stackrel{~}{s}_e\stackrel{~}{c}_e\stackrel{~}{v}_e}{(1+\stackrel{~}{v}_e^2)}(V(x)+V(y))}{(xM_Z^2)(yM_Z^2)}}]\}`$ (42)
$`G_2(x,y)=({\displaystyle \frac{3\mathrm{\Gamma }_e}{\alpha M_Z}}){\displaystyle \frac{1}{(1+\stackrel{~}{v}_e^2)}}[{\displaystyle \frac{\stackrel{~}{\mathrm{\Delta }}_\alpha (x)R(y)}{x(yM_Z^2)}}+{\displaystyle \frac{\stackrel{~}{\mathrm{\Delta }}_\alpha (y)R(x)}{y(xM_Z^2)}}]`$ (43)
$`({\displaystyle \frac{3\mathrm{\Gamma }_e}{\alpha M_Z}})^2{\displaystyle \frac{4\stackrel{~}{v}_e^2}{(1+\stackrel{~}{v}_e^2)^2}}[{\displaystyle \frac{R(x)+R(y)+\frac{4\stackrel{~}{s}_e\stackrel{~}{c}_e}{\stackrel{~}{v}_e}(V(x)+V(y))}{(xM_Z^2)(yM_Z^2)}}]\}`$ (44)
$`G_3(x,y)={\displaystyle \frac{\stackrel{~}{v}_e}{(1+\stackrel{~}{v}_e^2)}}\{({\displaystyle \frac{3\mathrm{\Gamma }_e}{\alpha M_Z}}){\displaystyle \frac{\stackrel{~}{\mathrm{\Delta }}_\alpha (x)R(y)\frac{4\stackrel{~}{s}_e\stackrel{~}{c}_e}{\stackrel{~}{v}_e}V(y)}{x(yM_Z^2)}}`$ (45)
$`({\displaystyle \frac{3\mathrm{\Gamma }_e}{\alpha M_Z}})^2[{\displaystyle \frac{R(x)+R(y)+\frac{4\stackrel{~}{s}_e\stackrel{~}{c}_e}{\stackrel{~}{v}_e}V(x)+\frac{8\stackrel{~}{s}_e\stackrel{~}{c}_e\stackrel{~}{v}_e}{(1+\stackrel{~}{v}_e^2)}V(y)}{(xM_Z^2)(yM_Z^2)}}]\}`$ (46)
where we use a condensed notation (for $`x,y`$ corresponding to $`q^2`$ or $`t`$), $`\stackrel{~}{\mathrm{\Delta }}_\alpha (q^2)`$ meaning $`\stackrel{~}{\mathrm{\Delta }}_\alpha (q^2,\theta )`$ and $`\stackrel{~}{\mathrm{\Delta }}_\alpha (t)`$ meaning $`\overline{\stackrel{~}{\mathrm{\Delta }}_\alpha }(q^2,\theta )`$; similarly for $`R`$ and $`V`$.
¿From the three previous quantities $`d\sigma ^{1,2,P}`$ we can compute for instance the unpolarized angular distribution
$$\frac{d\sigma }{dcos\theta }\frac{d\sigma ^1}{dcos\theta }+\frac{d\sigma ^2}{dcos\theta }$$
(47)
the Left-Right polarization asymmetry
$$A_{LR}(q^2,\theta )=[\frac{d\sigma ^P}{dcos\theta }]/[\frac{d\sigma }{dcos\theta }]$$
(48)
and the new (LL+RR)/(LR+RL+LL+RR) polarization asymmetry which arises from the typical $`t`$-channel scattering amplitude
$$A_{||}(q^2,\theta )=[\frac{d\sigma ^2}{dcos\theta }]/[\frac{d\sigma }{dcos\theta }]$$
(49)
## III Applications to several NP models
### A Universal NP with a high scale
The previous representation eqs.(37)-(46) continues to be valid in the presence of new physics (NP) that does not add extra Lorentz structures to those of the SM. In this case, one simply decomposes the three general one-loop functions as:
$$\stackrel{~}{\mathrm{\Delta }}_\alpha ,R,V=(\stackrel{~}{\mathrm{\Delta }}_\alpha ,R,V)^{SM}+(\stackrel{~}{\mathrm{\Delta }}_\alpha ,R,V)^{NP}$$
(50)
and compute the (NP) effects on the various observables, once their contribution to $`(\stackrel{~}{\mathrm{\Delta }}_\alpha ,R,V)`$ is specified.
For a model of new physics that does not satisfy special simplicity requests, the calculation of virtual effects in the Bhabha scattering is affected by a proliferation of terms with respect to the annihilation process $`e^+e^{}f\overline{f}`$, ($`fe`$), as one sees immediately from inspection of eqs.(37)-(46). In fact, after $`\theta `$-integration, one will find in general a set of different functions of $`q^2`$ that correspond to each power of $`\theta `$ in the integrand. Each set arises from the six original functions $`\stackrel{~}{\mathrm{\Delta }}_\alpha ,\overline{\stackrel{~}{\mathrm{\Delta }}}_\alpha ,R,\overline{R},V,\overline{V}`$, which means to double the corresponding number of the case $`fe`$. Although this can be a purely computational problem, it obviously complicates the practical treatment for this process.
The situation shows a drastic change for those models of new physics that satisfy the requests of being, at the same time, universal, independent of the $`s`$, $`t`$ channels scattering angle (e.g. only contributing self-energies and/or vertices), and endowed with an intrinsic scale $`\mathrm{\Lambda }`$ ”sufficiently” larger than $`\sqrt{q^2}`$. In fact, in the $`Z`$-peak subtracted approach, one has by construction
$$\stackrel{~}{\mathrm{\Delta }}_\alpha (0,\theta )=\overline{\stackrel{~}{\mathrm{\Delta }}}_\alpha (0,\theta )=R(M_Z^2,\theta )=\overline{R}(M_Z^2,\theta )=V(M_Z^2,\theta )=\overline{V}(M_Z^2,\theta )=0$$
(51)
For Universal New Physics effects of the previous considered type one can then write the following parametrization:
$$R^{UNP}(z)=\frac{(zM_Z^2)}{M_Z^2}[\delta _Z]$$
(52)
$$V^{UNP}(z)=\frac{(zM_Z^2)}{M_Z^2}[\delta _s]$$
(53)
$$\stackrel{~}{\mathrm{\Delta }}_\alpha ^{UNP}(z)=\frac{z}{M_Z^2}[\delta _\gamma ]$$
(54)
where $`z=q^2,t`$.
The quantities $`\delta _{Z,s,\gamma }`$ will be in general unknown functions of $`z`$. For $`\mathrm{\Lambda }^2>>z`$ we can reasonably assume that the three functions are smooth. This means that they could be well approximated by the coefficient of the lowest power in a $`q^2`$ expansion that is $`\delta _i(0)`$ whenever $`\delta _i(0)0`$ (this will be the case in the two considered examples). In this case, the same three parameters will describe the NP effects both on the $`s`$ and on the $`t`$ channel observables. These parameters are also the same that appear, for the chosen models, in all the remaining proccesses $`e^+e^{}f\overline{f}`$, ($`fe`$). This fact allows to combine the theoretical analysis of the two types of processes without increase of parameters, thus improving the accuracy of the conclusions that are reached.
The NP expression of the functions $`G_i(x,y)`$ acquires in this case the simple form:
$`G_1^{UNP}(x,y)={\displaystyle \frac{1}{M_Z^2}}\{\delta _\gamma [{\displaystyle \frac{1}{x}}+{\displaystyle \frac{1}{y}}]`$ (55)
$`+({\displaystyle \frac{3\mathrm{\Gamma }_e}{\alpha M_Z}})({\displaystyle \frac{\stackrel{~}{v}_e^2}{1+\stackrel{~}{v}_e^2}})[\delta _\gamma ({\displaystyle \frac{1}{xM_Z^2}}+{\displaystyle \frac{1}{yM_Z^2}})(\delta _Z+{\displaystyle \frac{8\stackrel{~}{s}_e\stackrel{~}{c}_e}{\stackrel{~}{v}_e}}\delta _s)({\displaystyle \frac{1}{x}}+{\displaystyle \frac{1}{y}})]`$ (56)
$`({\displaystyle \frac{3\mathrm{\Gamma }_e}{\alpha M_Z}})^2(\delta _Z+({\displaystyle \frac{8\stackrel{~}{s}_e\stackrel{~}{c}_e\stackrel{~}{v}_e}{1+\stackrel{~}{v}_e^2}})\delta _s)[{\displaystyle \frac{1}{xM_Z^2}}+{\displaystyle \frac{1}{yM_Z^2}}]\}`$ (57)
$`G_2^{UNP}(x,y)={\displaystyle \frac{1}{M_Z^2}}\{({\displaystyle \frac{3\mathrm{\Gamma }_e}{\alpha M_Z}})({\displaystyle \frac{1}{1+\stackrel{~}{v}_e^2}})[\delta _\gamma [{\displaystyle \frac{1}{xM_Z^2}}+{\displaystyle \frac{1}{yM_Z^2}}]\delta _Z[{\displaystyle \frac{1}{x}}+{\displaystyle \frac{1}{y}}]`$ (58)
$`({\displaystyle \frac{3\mathrm{\Gamma }_e}{\alpha M_Z}})^2{\displaystyle \frac{4\stackrel{~}{v}_e^2}{(1+\stackrel{~}{v}_e^2)^2}}(\delta _Z+{\displaystyle \frac{4\stackrel{~}{s}_e\stackrel{~}{c}_e}{\stackrel{~}{v}_e}}\delta _s)[{\displaystyle \frac{1}{xM_Z^2}}+{\displaystyle \frac{1}{yM_Z^2}}]\}`$ (59)
$`G_3^{UNP}(x,y)={\displaystyle \frac{1}{M_Z^2}}{\displaystyle \frac{\stackrel{~}{v}_e}{1+\stackrel{~}{v}_e^2}}\{({\displaystyle \frac{3\mathrm{\Gamma }_e}{\alpha M_Z}})[{\displaystyle \frac{\delta _\gamma }{yM_Z^2}}{\displaystyle \frac{\delta _Z+\frac{4\stackrel{~}{s}_e\stackrel{~}{c}_e}{\stackrel{~}{v}_e}\delta _s}{x}}]`$ (60)
$`({\displaystyle \frac{3\mathrm{\Gamma }_e}{\alpha M_Z}})^2[{\displaystyle \frac{\delta _Z+\frac{4\stackrel{~}{s}_e\stackrel{~}{c}_e}{\stackrel{~}{v}_e}\delta _s}{yM_Z^2}}+{\displaystyle \frac{\delta _Z+\frac{8\stackrel{~}{s}_e\stackrel{~}{c}_e\stackrel{~}{v}_e}{(1+\stackrel{~}{v}_e^2)}\delta _s}{xM_Z^2}}]\}`$ (61)
The three constants $`\delta _Z`$, $`\delta _s`$, $`\delta _\gamma `$ depend on the chosen model and can be easily determined in each separate case. To show how this procedure works in practice, we shall provide the expressions of the $`\delta _i`$ in a couple of specific cases that meet our simplicity requests. With this aim, we have considered the following models:
(1) Anomalous Gauge Couplings (AGC)
We used the framework of Ref. in which the effective Lagrangian is constructed with dimension six operators respecting $`SU(2)\times U(1)`$ and $`CP`$ invariance. As shown in Ref. , only two parameters ($`f_{DW}`$ and $`f_{DB}`$) survive in the Z-peak subtracted approach. The explicit expression of the UNP contribution to $`\delta _Z`$, $`\delta _s`$ and $`\delta _\gamma `$ are
$`\delta _Z=8\pi \alpha ({\displaystyle \frac{M_Z^2}{\mathrm{\Lambda }^2}})({\displaystyle \frac{\stackrel{~}{c}_e^2}{\stackrel{~}{s}_e^2}}f_{DW}+{\displaystyle \frac{\stackrel{~}{s}_e^2}{\stackrel{~}{c}_e^2}}f_{DB})`$ (62)
$`\delta _s=8\pi \alpha ({\displaystyle \frac{M_Z^2}{\mathrm{\Lambda }^2}})({\displaystyle \frac{\stackrel{~}{c}_e}{\stackrel{~}{s}_e}}f_{DW}{\displaystyle \frac{\stackrel{~}{s}_e}{\stackrel{~}{c}_e}}f_{DB})`$ (63)
$`\delta _\gamma =8\pi \alpha ({\displaystyle \frac{M_Z^2}{\mathrm{\Lambda }^2}})(f_{DW}+f_{DB})`$ (64)
They satisfy the linear constraint:
$$\delta _Z\frac{12\stackrel{~}{s}_e^2}{\stackrel{~}{s}_e\stackrel{~}{c}_e}\delta _s+\delta _\gamma =0.$$
(65)
(2) Technicolour
The second considered model was one of Technicolour type with two families of strongly coupled resonances ($`V`$ and $`A`$) . The typical UNP parameters are the two ratios $`F_A/M_A`$ and $`F_V/M_V`$ where $`F_{A,V}`$ and $`M_{A,V}`$ are the couplings and the masses (that in this case play the role of the new physics scale $`\mathrm{\Lambda }^{TC}>>q^2`$) of the lightest axial and vector resonances. The contribution to $`\delta _Z`$, $`\delta _s`$ and $`\delta _\gamma `$ are
$`\delta _Z={\displaystyle \frac{\pi \alpha }{\stackrel{~}{s}_e^2\stackrel{~}{c}_e^2}}[(12\stackrel{~}{s}_e^2)^2{\displaystyle \frac{M_Z^2}{M_V^4}}F_V^2+{\displaystyle \frac{M_Z^2}{M_A^4}}F_A^2]`$ (66)
$`\delta _s={\displaystyle \frac{2\pi \alpha }{\stackrel{~}{s}_e\stackrel{~}{c}_e}}(12\stackrel{~}{s}_e^2)(12\stackrel{~}{s}_e^2)^2{\displaystyle \frac{M_Z^2}{M_V^4}}F_V^2`$ (67)
$`\delta _\gamma =4\pi \alpha ({\displaystyle \frac{M_Z^2}{M_V^4}})F_V^2`$ (68)
Again, we have a linear constraint in the $`(\delta _Z,\delta _s,\delta _\gamma )`$ space:
$$\delta _s=\left(\frac{12\stackrel{~}{s}_e^2}{2\stackrel{~}{s}_e\stackrel{~}{c}_e}\right)\delta _\gamma .$$
(69)
and the conditions
$$\delta _{Z,s}>0\delta _\gamma <0$$
(70)
### B Non universal examples
Strictly speaking, our procedure has been motivated by the possibility of investigating models of a special universal type, for which the number of parameters to be determined can be suitably reduced. But there exist interesting models of NP that, although not of universal type, can be nevertheless described by a very restricted number of parameters. In these special simple cases our $`Z`$-peak subtracted procedure can be applied, without invoking any smoothness assumption, using the more general expressions of eqs.(37)-(46). In what follows, we have considered two cases that seem to us particularly relevant. These are:
(3) Contact terms
With the idea of compositeness (but it applies to any virtual NP effect with a high intrinsic scale, for example higher vector boson exchanges, satisfying chirality conservation) the following interaction
$``$ $`=`$ $`k_{if}{\displaystyle \frac{4\pi }{\mathrm{\Lambda }^2}}\{\eta _{LL}(\overline{\mathrm{\Psi }}_L^i\gamma ^\mu \mathrm{\Psi }_L^i)(\overline{\mathrm{\Psi }}_L^f\gamma _\mu \mathrm{\Psi }_L^f)+\eta _{RR}(\overline{\mathrm{\Psi }}_R^i\gamma ^\mu \mathrm{\Psi }_R^i)(\overline{\mathrm{\Psi }}_R^f\gamma _\mu \mathrm{\Psi }_R^f)`$ (72)
$`+\eta _{RL}(\overline{\mathrm{\Psi }}_R^i\gamma ^\mu \mathrm{\Psi }_R^i)(\overline{\mathrm{\Psi }}_L^f\gamma _\mu \mathrm{\Psi }_L^f)+\eta _{LR}(\overline{\mathrm{\Psi }}_L^i\gamma ^\mu \mathrm{\Psi }_L^i)(\overline{\mathrm{\Psi }}_R^f\gamma _\mu \mathrm{\Psi }_R^f)\}`$
was first introduced in for any four fermion interaction $`(i\overline{i}f\overline{f})`$; $`k_{if}=\frac{1}{2}`$ for $`if`$, $`k_{if}=1`$ otherwise; $`\mathrm{\Psi }_L=(1\gamma ^5)/2\mathrm{\Psi }`$, $`\mathrm{\Psi }_R=(1+\gamma ^5)/2\mathrm{\Psi }`$; $`\eta _{ab}`$ are phase factors defining the chirality structure of the interaction. Various applications have been made for pure chiral cases $`(ij)=LL`$ or $`RR`$ or $`LR`$ or $`RL`$ (keeping only one $`\eta _{ij}=\pm 1`$), as well as for mixed cases like $`VV`$ ($`\eta _{LL}=\eta _{RR}=\eta _{RL}=\eta _{LR}=\pm 1`$), $`AA`$ ($`\eta _{LL}=\eta _{RR}=\eta _{RL}=\eta _{LR}=\pm 1`$), $`VA`$ ($`\eta _{LL}=\eta _{RR}=\eta _{RL}=\eta _{LR}=\pm 1`$), $`AV`$ ($`\eta _{LL}=\eta _{RR}=\eta _{RL}=\eta _{LR}=\pm 1`$); see ref. for a general discussion.
In the $`Z`$-peak subtracted representation, the effect of this interaction on the $`e^+e^{}f\overline{f}`$ ($`fe`$) observables is obtained through the following expressions:
$`\stackrel{~}{\mathrm{\Delta }}_{\alpha ,ef}(q^2,\theta )=({\displaystyle \frac{\pi q^2}{e^2Q_eQ_f\mathrm{\Lambda }^2}})[\eta _{LL}(1v_e)(1v_f)+\eta _{RR}(1+v_e)(1+v_f)`$ (73)
$`+\eta _{RL}(1+v_e)(1v_f)+\eta _{LR}(1v_e)(1+v_f)]`$ (74)
$`R_{ef}(q^2,\theta )=({\displaystyle \frac{4\stackrel{~}{s}_e^2\stackrel{~}{c}_e^2\pi (q^2M_Z^2)}{e^2I_{3e}I_{3f}\mathrm{\Lambda }^2}})[\eta _{LL}+\eta _{RR}\eta _{RL}\eta _{LR}]`$ (75)
$`V_{ef}^{\gamma Z}(q^2,\theta )=({\displaystyle \frac{2\stackrel{~}{s}_e\stackrel{~}{c}_e\pi (q^2M_Z^2)}{e^2Q_eI_{3f}\mathrm{\Lambda }^2}})[\eta _{LL}(1v_e)\eta _{RR}(1+v_e)+\eta _{RL}(1+v_e)\eta _{LR}(1v_e)]`$ (76)
$`V_{ef}^{Z\gamma }(q^2,\theta )=({\displaystyle \frac{2\stackrel{~}{s}_e\stackrel{~}{c}_e\pi (q^2M_Z^2)}{e^2Q_fI_{3e}\mathrm{\Lambda }^2}})[\eta _{LL}(1v_f)\eta _{RR}(1+v_f)\eta _{RL}(1v_f)+\eta _{LR}(1+v_f)]`$ (77)
For each choice of chirality structure, like $`LL`$, $`RR`$, $`LR`$, $`RL`$, $`VV`$,$`AA`$,$`VA`$,$`AV`$, the effects on the differential cross section for two fermion production can be described by a single parameter $`\mathrm{\Lambda }`$.
In the case of Bhabha scattering, the constraint $`\eta _{RL}=\eta _{LR}`$ applies, so that the above expression can be put in the form of eq.(52-54), with:
$`\delta _Z=({\displaystyle \frac{16\stackrel{~}{s}_e^2\stackrel{~}{c}_e^2\pi M_Z^2}{e^2\mathrm{\Lambda }^2}})[\eta _{LL}+\eta _{RR}2\eta _{LR}]`$ (78)
$`\delta _s=({\displaystyle \frac{4\stackrel{~}{s}_e\stackrel{~}{c}_e\pi M_Z^2}{e^2\mathrm{\Lambda }^2}})[\eta _{LL}(1v_e)\eta _{RR}(1+v_e)+2v_e\eta _{LR}]`$ (79)
$`\delta _\gamma =({\displaystyle \frac{\pi M_Z^2}{e^2\mathrm{\Lambda }^2}})[\eta _{LL}(1v_e)^2+\eta _{RR}(1+v_e)^2+2\eta _{LR}(1v_e^2)]`$ (80)
(4) Extra dimensions
Recently, an intense activity has been developed on possible low energy effects of graviton exchange. The following matrix element for the 4-fermion process $`e^+e^{}\overline{f}f`$ is predicted:
$$\frac{\lambda }{\mathrm{\Lambda }^4}[\overline{e}\gamma ^\mu e\overline{f}\gamma _\mu f(p_2p_1).(p_4p_3)\overline{e}\gamma ^\mu e\overline{f}\gamma ^\nu f(p_2p_1)_\nu (p_4p_3)_\mu ]$$
(81)
For this model one finds in the case of Bhabha scattering the following contributions to $`\delta _{\gamma ,Z,s}`$ that, as one sees, are now genuine $`q^2`$, $`\theta `$ functions:
$`\delta _\gamma =({\displaystyle \frac{\lambda q^2M_Z^2}{\mathrm{\Lambda }^4}}){\displaystyle \frac{(\stackrel{~}{v}_e^22cos\theta )}{e^2}}`$ (82)
$`\delta _Z=({\displaystyle \frac{\lambda q^2M_Z^2}{\mathrm{\Lambda }^4}})({\displaystyle \frac{16\stackrel{~}{s}_e^2\stackrel{~}{c}_e^2}{e^2}})`$ (83)
$`\delta _s=({\displaystyle \frac{\lambda q^2M_Z^2}{\mathrm{\Lambda }^4}})({\displaystyle \frac{4\stackrel{~}{s}_e\stackrel{~}{c}_e\stackrel{~}{v}_e}{e^2}})`$ (84)
Illustrations will be given in the next Section with the normalization $`\lambda =\pm 1`$. Note that the $`q^2`$ factor is purely kinematical and a consequence of the higher dimension of the interaction Lagrangian. Note also the presence of the term proportional to $`q^2cos\theta `$ in the $`s`$-channel photon coefficient $`\delta _\gamma `$, which gives a contribution proportional to $`t+2q^2`$ in the $`t`$-channel, according to eq.(26). This contribution will turn out to give the largest effect through the interference with the standard photon exchange amplitude.
Our theoretical description of new physics effects is at this point concluded. The final Section 4 will be devoted to a detailed numerical analysis of the information that can be derived on the involved parameters by the present (LEP2) and future (LC), using both Bhabha scattering and all the remaining $`e^+e^{}f\overline{f}`$ processes.
## IV Numerical Results
### A LEP2 (present and future)
As a first application of our approach, we have used some of the LEP2 results on fermion pair production in order to constrain the set of constants $`\delta _{Z,s,\gamma }`$ which fully describes the effects of general Universal New Physics. In particular, in the same spirit of a previous study , we have considered the following “non Bhabha” observables: $`\sigma _{\mu ,\tau }`$ (the cross section for $`\mu `$ and $`\tau `$ pair production), $`A_{FB,\mu ,\tau }`$ (the related forward-backward asymmetries) and $`\sigma _5`$ (the cross section for production of quark pairs, for the five light flavours accessible at LEP energies). In addition to these observables we have included the unpolarized differential Bhabha cross section, measured in intervals of the cosine of the polar angle of the scattered electron. For the muon and tau cross sections and asymmetries and for the hadronic cross section the combinations of preliminary results of the four LEP experiments with the data collected at center of mass energies of 183 GeV and 189 GeV have been used . Although based on preliminary results, the combined measurements allow to take advantage of the whole data sample produced at LEP2 and therefore to benefit of the reduced statistical error and of the proper treatment of the various sources of experimental systematic uncertainty. A measurement of the differential cross section for Bhabha scattering with acollinearity smaller that 10 has been recently performed with a data sample of approximately 180 $`pb^1`$ at the center of mass energy of 189 GeV . The differential cross section is measured in nine uniform intervals of the polar angle of the scattered electron, $`\mathrm{cos}\theta _e^{}`$, in the range (-0.9, 0.9). The precision of the measurement, which is limited primarely by the statistical uncertainty, reaches the level of 1% in the interval of most forward scattering angles.
This measurement, together with the combined results on muon, tau and hadronic observables, has been compared to the Standard Model prediction corresponding to the experimental signal definition. The deviations of the measurement with respect to the Standard Model expectations have then been used to measure and constrain the parameters $`\delta _Z`$, $`\delta _s`$ and $`\delta _\gamma `$ with a $`\chi ^2`$ fit. In the fit procedure the uncertainty on the reference Standard Model prediction itself must be taken into account. The theoretical uncertainties on the Standard Model predictions for fermion pair production at LEP2 energies are mainly related to the estimate of the large QED corrections. In the case of $`\mu ^+\mu ^{}`$, $`\tau ^+\tau ^{}`$ and $`q\overline{q}`$ production, the differences between predictions of several semianalytic or Monte Carlo calculations for cross sections and asymmetries are smaller than 1% and, therefore, they are negligible with respect to the experimental error. In the Bhabha scattering process the different programs providing the Standard Model predictions compare each other at the level of 2% in the experimental acceptance. Therefore, in our study, we have assigned a 2% uncertainty to the reference Standard Model prediction for the Bhabha differential cross section. This uncertainty reflects into an error larger than the experimental one in the region of forward scattering, which, as will be discussed in the following, is the most sensitive to New Physics effects in $`\delta _\gamma `$.
The results of the analysis are shown in Tab. (I). As one sees, the addition of Bhabha scattering improves, although not spectacularly, the bound on $`\delta _\gamma `$ which is constrained by the data in the forward scattering angle region, where the Bhabha cross section is dominated by the $`t`$-channel photon exchange contribution. It is interesting to take in mind that , at present, the sensitivity to NP effects in the forward Bhabha cross section is spoiled by the theoretical uncertainty on the Standard Model theoretical prediction, which dominates over the experimental error. Should this error be reduced, the role of this measurements would certainly be more relevant.
To give a more quantitative meaning to the latter claim, we have simulated a forthcoming measurement at 200 GeV with an overall $`400`$ $`pb^1`$ luminosity for each experiment, and repeated the previous analysis adding these future data to those available at $`183`$, $`189`$ GeV. For consistency, we have assumed in all the three sets of data a central value coincident with the SM prediction. The errors are those available at $`183`$ and $`189`$ GeV. At $`200`$ GeV we considered two scenarios, one with purely statistical errors and one with their combination with a 2% theoretical error on Bhabha scattering. The difference between the two cases affects mainly the very forward scattering cone and therefore $`\delta _\gamma `$. The results of this second analysis are shown in Tab. (II). As one sees, in agreement with the qualitative expectations, the role of Bhabha scattering is now definitely more relevant in the determination of the bound for the photonic parameter $`\delta _\gamma `$.
The same two analyses have been performed for the two models involving contact terms and extra dimensions. The results are presented in Tables III,IV. For what concerns the contact terms, one should first note that the values of the bounds obtained from $`e^+e^{}\mu ^+\mu ^{},q\overline{q}`$ (without Bhabha) depend strongly on the chirality structure. This comes from the interference of the contact amplitude with the standard $`\gamma ,Z`$ exchange amplitude, which is larger when the chirality structures of both amplitudes are close to each other. In particular the $`VA`$ bound is found very low because the standard $`VA`$ amplitude is depressed by the small $`Ze^+e^{}`$ vector coupling. One then sees that the effect of the Bhabha process on these bounds is generally modest. An opposite situation appears for the case of extra dimensions, where the bulk of the effect is provided by the “forward” data. As already mentioned in Section IIIB(4), the largest effect on the differential cross section comes from the interference of the standard photon exchange with the extra dimension term both in the $`t`$-channel, followed by the mixed ($`s`$-channel)- ($`t`$-channel) one, the ($`s`$-channel)-($`s`$-channel) contribution being much smaller. Once again, this is more clearly visible in the analysis that uses the future data at 200 GeV, in the case one assumes no theoretical error in the the Bhabha component.
All the numerical results exhibited in Tabs. (I,II) can be represented graphically. We have shown in Figs. (1,2) the planar ellipses that are obtained by projecting onto the three planes $`(\delta _s,\delta _\gamma )`$, $`(\delta _Z,\delta _\gamma )`$ and $`(\delta _Z,\delta _s)`$ the 95 % C. L. allowed three dimensional region resulting from a global fit of all data in terms of the three parameters $`\delta _{Z,s,\gamma }`$. For completeness, we also show the results for the two representative AGC and TC models. As a very preliminary comment concerning the latter cases we can notice that, although at a rather qualitative level, LEP2 data apparently do not particularly support the considered TC theoretical proposal that would require $`\delta _s>0`$.
### B LC analyses
This analysis has been performed in a spirit that is very similar to that used for the future 200 GeV LEP2 analyses. In other words, we have assumed a set of measurements at $`\sqrt{q^2}=500`$ GeV whose central values agree with the SM predictions, and postulated a purely statistical error corresponding to a high luminosity of 500 $`fb^1`$. We have added to the previous LEP2 “non Bhabha” observables the longitudinal polarization asymmetry $`A_{LR}`$ for lepton production. A discussion of the important role of this observable can be found in . Of course, in principle other measurements e.g. for final $`b`$ or $`t`$ quarks could be used. We have also assumed 9 angular Bhabha measurements for all the three different observables ($`\sigma `$, $`\sigma ^1`$, $`\sigma ^P`$) defined in Section (2). More precisely, in each bin $`[\theta _{\mathrm{min}},\theta _{max}]`$ we have considered
$$\sigma _{[\theta _{min},\theta _{max}]}=_{\mathrm{cos}\theta _{max}}^{\mathrm{cos}\theta _{min}}\frac{d\sigma }{d\mathrm{cos}\theta }d\mathrm{cos}\theta $$
(85)
$$A_{LR,[\theta _{min},\theta _{max}]}=\frac{1}{\sigma _{[\theta _{min},\theta _{max}]}}_{\mathrm{cos}\theta _{max}}^{\mathrm{cos}\theta _{min}}\frac{d\sigma ^P}{d\mathrm{cos}\theta }d\mathrm{cos}\theta ,$$
(86)
$$A_{||,[\theta _{min},\theta _{max}]}=\frac{1}{\sigma _{[\theta _{min},\theta _{max}]}}_{\mathrm{cos}\theta _{max}}^{\mathrm{cos}\theta _{min}}\frac{d\sigma ^2}{d\mathrm{cos}\theta }d\mathrm{cos}\theta ,$$
(87)
Tables (V, VI) contain the numerical results for the universal and not universal models. As a general feature, one notices that in the universal case the limits on all the three parameters $`\delta _{Z,s,\gamma }`$ are substantially (a factor of two) improved by the use of Bhabha observables. The interesting feature is that, in each case, different Bhabha observables play the crucial role. In fact, $`\delta _Z`$ is mostly affected by $`A_{||}`$ (in both angular directions), $`\delta _s`$ (as one expects) by $`A_{LR}`$ ( in the forward cone), and $`\delta _\gamma `$ by the unpolarized $`\sigma `$ (again, for very small angles). In the considered non universal cases, the effect is, again, not spectacular (although not negligible) for the contact terms. Quite on the contrary, there would be a large (a factor two) effect in the case of extra dimensions, mostly due to the unpolarized cross section at small angles. For this specific model of new physics Bhabha scattering seems therefore to represent a fundamental experimental measurement. This statement is well in agreement with the results of a recent numerical analysis of LEP2 data .
## V Conclusion
In conclusion, we have shown that, for a class of theoretical models of new physics that is certainly not empty, the generalization of the $`Z`$-peak subtracted approach to the case of Bhabha scattering can be simply performed, leading in general to improvements of the information that might be obtained. As a general statement, Bhabha scattering appears to be always relevant; for models of universal type, polarized and unpolarized observables play a crucial role in the determination of the bounds for the different parameters; in other non universal interesting cases, like in particular that of extra dimensions, unpolarized Bhabha observables appear to play a fundamental role.
Appendix A: General form of the polarized Bhabha scattering cross section.
$$\frac{d\sigma }{dcos\theta }=(1PP^{})\frac{d\sigma ^1}{dcos\theta }+(1+PP^{})\frac{d\sigma ^2}{dcos\theta }+(P^{}P)\frac{d\sigma ^P}{dcos\theta }$$
(88)
with
$`{\displaystyle \frac{d\sigma ^{N1}}{dcos\theta }}={\displaystyle \frac{1}{16\pi q^2}}\{{\displaystyle \frac{t^2+u^2}{q^4}}(g_{Ve}^{(\gamma )})^4+{\displaystyle \frac{2(q^2M_Z^2)}{q^2[(q^2M_Z^2)^2+M_Z^2\mathrm{\Gamma }_Z^2]}}[(t^2+u^2)(g_{Ve}^{(\gamma )}g_{Ve}^{(Z)})^2+(u^2t^2)(g_{Ve}^{(\gamma )}g_{Ae}^{(Z)})^2]`$ (89)
$`+{\displaystyle \frac{1}{[(q^2M_Z^2)^2+M_Z^2\mathrm{\Gamma }_Z^2]}}[(t^2+u^2)[(g_{Ve}^{(Z)})^2+(g_{Ae}^{(Z)})^2]+4(u^2t^2)(g_{Ve}^{(Z)}g_{Ae}^{(Z)})^2]+{\displaystyle \frac{2u^2}{q^2t}}(g_{Ve}^{(\gamma )}\overline{g}_{Ve}^{(\gamma )})^2`$ (90)
$`+{\displaystyle \frac{2u^2}{q^2(tM_Z^2)}}(g_{Ve}^{(\gamma )})^2[(\overline{g}_{Ve}^{(Z)})^2+(\overline{g}_{Ae}^{(Z)})^2]+{\displaystyle \frac{2u^2(q^2M_Z^2)}{t[(q^2M_Z^2)^2+M_Z^2\mathrm{\Gamma }_Z^2]}}(\overline{g}_{Ve}^{(\gamma )})^2[(g_{Ve}^{(Z)})^2+(g_{Ae}^{(Z)})^2]`$ (91)
$`+{\displaystyle \frac{2u^2(q^2M_Z^2)}{(tM_Z^2)[(q^2M_Z^2)^2+M_Z^2\mathrm{\Gamma }_Z^2]}}([(g_{Ve}^{(Z)})^2+(g_{Ae}^{(Z)})^2][(\overline{g}_{Ve}^{(Z)})^2+(\overline{g}_{Ae}^{(Z)})^2]+4g_{Ve}^{(Z)}g_{Ae}^{(Z)}\overline{g}_{Ve}^{(Z)}\overline{g}_{Ae}^{(Z)})`$ (92)
$`+{\displaystyle \frac{u^2}{t^2}}(\overline{g}_{Ve}^{(\gamma )})^4+{\displaystyle \frac{2u^2}{t(tM_Z^2)}}(\overline{g}_{Ve}^{(\gamma )})^2[(\overline{g}_{Ve}^{(Z)})^2+(\overline{g}_{Ae}^{(Z)})^2]+{\displaystyle \frac{u^2}{(tM_Z^2)^2}}[((\overline{g}_{Ve}^{(Z)})^2+(\overline{g}_{Ae}^{(Z)})^2)^2+4(\overline{g}_{Ve}^{(Z)}\overline{g}_{Ae}^{(Z)})^2]\}`$ (93)
(94)
$`{\displaystyle \frac{d\sigma ^2}{dcos\theta }}={\displaystyle \frac{1}{16\pi q^2}}\{{\displaystyle \frac{q^4}{t^2}}(\overline{g}_{Ve}^{(\gamma )})^4+{\displaystyle \frac{2q^4}{t(tM_Z^2)}}(\overline{g}_{Ve}^{(\gamma )})^2[(g_{Ve}^{(Z)})^2(g_{Ae}^{(Z)})^2]`$ (95)
$`+{\displaystyle \frac{q^4}{(tM_Z^2)^2}}([(\overline{g}_{Ve}^{(Z)})^2(\overline{g}_{Ae}^{(Z)})^2]^2)\}`$ (96)
$`{\displaystyle \frac{d\sigma ^P}{dcos\theta }}={\displaystyle \frac{u^2}{4\pi q^2}}\{{\displaystyle \frac{1}{[(q^2M_Z^2)^2+M_Z^2\mathrm{\Gamma }_Z^2]}}({\displaystyle \frac{q^2M_Z^2}{q^2}}(g_{Ve}^{(\gamma )})^2(g_{Ve}^{(Z)}g_{Ae}^{(Z)})+(g_{Ve}^{(Z)}g_{Ae}^{(Z)})[(g_{Ve}^{(Z)})^2+(g_{Ae}^{(Z)})^2])`$ (97)
$`+{\displaystyle \frac{1}{q^2(tM_Z^2)}}(g_{Ve}^{(\gamma )})^2(\overline{g}_{Ve}^{(Z)}\overline{g}_{Ae}^{(Z)})+{\displaystyle \frac{q^2M_Z^2)}{t[(q^2M_Z^2)^2+M_Z^2\mathrm{\Gamma }_Z^2]}}(\overline{g}_{Ve}^{(\gamma )})^2(g_{Ve}^{(Z)}g_{Ae}^{(Z)})`$ (98)
$`+{\displaystyle \frac{q^2M_Z^2}{(tM_Z^2)[(q^2M_Z^2)^2+M_Z^2\mathrm{\Gamma }_Z^2]}}((g_{Ve}^{(Z)}g_{Ae}^{(Z)})[(\overline{g}_{Ve}^{(Z)})^2+(\overline{g}_{Ae}^{(Z)})^2]+(\overline{g}_{Ve}^{(Z)}\overline{g}_{Ae}^{(Z)})[(g_{Ve}^{(Z)})^2+(g_{Ae}^{(Z)})^2])`$ (99)
$`+{\displaystyle \frac{1}{t(tM_Z^2)}}(\overline{g}_{Ve}^{(\gamma )})^2(\overline{g}_{Ve}^{(Z)}\overline{g}_{Ae}^{(Z)})+{\displaystyle \frac{1}{(tM_Z^2)^2}}(\overline{g}_{Ve}^{(Z)}\overline{g}_{Ae}^{(Z)}[(\overline{g}_{Ve}^{(Z)})^2+(\overline{g}_{Ae}^{(Z)})^2]\}`$ (100) |
warning/0002/astro-ph0002249.html | ar5iv | text | # 1 Abstract
## 1 Abstract
I review the basic theory of the cosmic microwave background (CMB) anisotropies in adiabatic cold dark matter (CDM) cosmologies. The latest observational results on the CMB power spectrum are consistent with the simplest inflationary models and indicate that the Universe is close to spatially flat with a nearly scale invariant fluctuation spectrum. We are also beginning to see interesting constraints on the density of CDM, with a best fit value of $`\omega _c\mathrm{\Omega }_ch^20.1`$. The CMB constraints, when combined with observations of distant Type Ia supernovae, are converging on a $`\mathrm{\Lambda }`$-dominated Universe with $`\mathrm{\Omega }_m0.3`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }0.7`$.<sup>1</sup><sup>1</sup>1To appear in Proceedings of NATO ASI: Structure formation in the Universe, eds. N. Turok, R. Crittenden.
## 2 Introduction
The discovery of temperature anisotropies in the CMB by the COBE team (Smoot et al. 1992) heralded a new era in cosmology. For the first time COBE provided a clear detection of the primordial fluctuations responsible for the formation of structure in the Universe at a time when they were still in the linear growth regime. Since then, a large number of ground based and balloon borne experiments have been performed which have succeeded in defining the shape of the power spectrum of temperature anisotropies $`C_{\mathrm{}}`$<sup>2</sup><sup>2</sup>2The power spectrum is defined as $`C_{\mathrm{}}=|a_\mathrm{}m|^2`$, where the $`a_\mathrm{}m`$ are determined from a spherical harmonic expansion of the temperature anisotropies on the sky, $`\mathrm{\Delta }T/T=a_\mathrm{}mY_\mathrm{}m(\theta ,\varphi )`$. up to multipoles of $`\mathrm{}300`$ clearly defining the first acoustic peak in the spectrum. Figure 1 shows a compilation of band power anisotropy measurements
$$\frac{\mathrm{\Delta }T_{\mathrm{}}}{T}=\sqrt{\frac{1}{2\pi }\mathrm{}(\mathrm{}+1)C_{\mathrm{}}}$$
(1)
that is almost up to date at the time of writing. The horizontal error bars show the multipole range probed by each experiment. The recent results from the VIPER experiment (Peterson et al. 1999) and the Boomerang test flight (Mauskopf et al. 1999) are not plotted because the exact window functions are not yet publically available. Neither are the published results from the Python V experiment (Coble et al. 1999) which seem to be discrepant with the other experiments particularly in the multipole range $`\mathrm{}<100`$. The points plotted in figure 1 are generally consistent with each other and provide strong evidence for a peak in the power spectrum at $`\mathrm{}200`$.
In this introductory article, I will review briefly the theory of CMB anisotropies in adiabatic models of structure formation and then discuss the implications of Figure 1 for values of cosmological parameters. The literature on the CMB anisotropies has grown enormously over the last few years and it is impossible to do the subject justice in a short article. General reviews of the CMB anisotropies are given by Bond (1996) and Kamionkowski and Kosowsky (1999). A recent review on constraining cosmological parameters from the CMB is given by Rocha (1999).
## 3 Basic Theory
Most of the key features of figure 1 can be understood using a simplified set of equations. The background universe is assumed to be spatially flat together with small perturbations $`h_{ij}`$ so that the metric is
$`ds^2=a^2(\tau )\left(\eta _{ij}+h_{ij}\right)dx^idx^j,`$ (2)
$`\eta _{ij}=(1,1,1,1),\tau ={\displaystyle 𝑑t/a}.`$
We adopt the synchronous gauge, $`h_{00}=h_{0i}=0`$, and ignore the anisotropy of Thomson scattering and perturbations in the relativistic neutrino component. With these assumptions, the equations governing the evolution of scalar plane wave perturbations of wavenumber $`k`$ are
$`\dot{\mathrm{\Delta }}+ik\mu \mathrm{\Delta }+\mathrm{\Phi }=\sigma _Tn_ea\left[\mathrm{\Delta }_0+4\mu v_b\mathrm{\Delta }\right]`$ (3a)
$`\mathrm{\Phi }=(3\mu ^21)\dot{h}_{33}(1\mu ^2)\dot{h}`$ (3b)
$`\dot{v}_b+{\displaystyle \frac{\dot{a}}{a}}v_b=\sigma _Tn_ea{\displaystyle \frac{\overline{\rho }_\gamma }{\overline{\rho }_b}}\left(\mathrm{\Delta }_1{\displaystyle \frac{4}{3}}v_b\right),`$ (3c)
$`\dot{\delta }_b={\displaystyle \frac{1}{2}}\dot{h}ikv_b,\dot{\delta }_C={\displaystyle \frac{1}{2}}\dot{h}`$ (3d)
$`\ddot{h}+{\displaystyle \frac{\dot{a}}{a}}\dot{h}=8\pi Ga^2(\overline{\rho }_b\delta _b+\overline{\rho }_c\delta _c+2\overline{\rho }_\gamma \mathrm{\Delta }_0)`$ (3e)
$`ik(\dot{h}_{33}\dot{h})=16\pi Ga^2(\overline{\rho }_bv_b+\overline{\rho }_\gamma \mathrm{\Delta }_1).`$ (3f)
Here, $`\mathrm{\Delta }`$ is the perturbation to the photon radiation brightness and $`\mathrm{\Delta }_0`$ and $`\mathrm{\Delta }_1`$ are its zeroth and first angular moments, $`\delta _b`$ and $`\delta _c`$ are the density perturbations in the baryonic and CDM components, $`v_b`$ is the baryon velocity and $`h=\mathrm{Tr}(h_{ij})`$. Dots denote differentiation with respect to the conformal time variable $`\tau `$. It is instructive to look at the solutions to these equations in the limits of large ($`k\tau _R1`$) and small ($`k\tau _R1`$) perturbations, where $`\tau _R`$ is the conformal time at recombination:
### 3.1 Large angle anisotropies
In the limit $`k\tau _R1`$, Thomson scattering is unimportant and so the term in square brackets in the Boltzmann equation for the photons can be ignored. In the matter dominated era $`h_{33}=h\tau ^2`$ and so equation (3a) becomes
$$\dot{\mathrm{\Delta }}+ik\mu \mathrm{\Delta }=2\mu ^2\dot{h}$$
(4)
with approximate solution
$$\mathrm{\Delta }(k,\mu ,\tau )\frac{2\ddot{h}(\tau _R)}{k^2}\mathrm{exp}\left(ik\mu (\tau _s\tau )\right).$$
(5)
This solution is the Sachs-Wolfe (1967) effect. Any deviation from the evolution $`\ddot{h}=\mathrm{constant}`$, caused for example by a non-zero cosmological constant, will lead to additional terms in equation (5) increasing the large-angle anisotropies (sometimes referred to as the late-time Sachs-Wolfe effect, see e.g. Bond 1996). The CMB power spectrum is given by
$$C_{\mathrm{}}=\frac{1}{8\pi }_0^{\mathrm{}}|\mathrm{}_{\mathrm{}}|^2k^2𝑑k,$$
(6)
where the perturbation $`\mathrm{\Delta }`$ has been expanded in Legendre polynomials,
$$\mathrm{\Delta }=\underset{\mathrm{}}{}(2\mathrm{}+1)\mathrm{\Delta }_{\mathrm{}}P_{\mathrm{}}(\mu ).$$
(7)
Inserting the solution of equation (5) into equation (6) gives
$$C_{\mathrm{}}=\frac{1}{2\pi }_0^{\mathrm{}}\frac{|\ddot{h}|^2}{k^4}j_l^2(k\tau _0)k^2𝑑k,$$
(8)
and so for a power-law spectrum of scalar perturbations, $`|h|^2k^{n_s}`$, the CMB power spectrum is
$$C_{\mathrm{}}=C_2\frac{\mathrm{\Gamma }\left(\mathrm{}+\frac{(n_s1)}{2}\right)}{\mathrm{\Gamma }\left(\mathrm{}+\frac{(5n_s)}{2}\right)}\frac{\mathrm{\Gamma }\left(\frac{9n_s}{2}\right)}{\mathrm{\Gamma }\left(\frac{3+n_s}{2}\right)}$$
(9)
giving the characteristic power-law like form, $`C_{\mathrm{}}\mathrm{}^{(n_s3)}`$ at low multipoles ($`\mathrm{}<30`$).
### 3.2 Small angle anisotropies and Acoustic peaks
In the matter dominated era, equation (3a) becomes
$$\dot{\mathrm{\Delta }}+ik\mu \mathrm{}=\sigma _Tn_ea\left[\mathrm{\Delta }_0+4\mu v_b\mathrm{\Delta }\right]+2\mu ^2\dot{h},$$
(10)
and taking the zeroth and first angular moments gives
$`\dot{\mathrm{\Delta }}_0+ik\mathrm{\Delta }_1={\displaystyle \frac{2}{3}}\dot{h}`$ (11a)
$`\dot{\mathrm{\Delta }}_1+ik\left({\displaystyle \frac{\mathrm{\Delta }_0+2\mathrm{\Delta }_2}{3}}\right)=\sigma _Tn_ea\left[{\displaystyle \frac{4}{3}}v_b\mathrm{\Delta }_1\right].`$ (11b)
Prior to recombination, $`\tau /\tau _c1`$ where $`\tau _c=1/(\sigma _Tn_ea)`$ is the mean collision time, and so the matter is tightly coupled to the radiation. In this limit $`\mathrm{\Delta }_14/3v_b`$ from equation (3c) and $`\mathrm{\Delta }_2`$ in equation (11b) can be ignored. With these approximations, equation (11b) becomes
$`\dot{\mathrm{\Delta }}_1+{\displaystyle \frac{ik\mathrm{\Delta }_0}{3}}={\displaystyle \frac{\overline{\rho }_b}{\overline{\rho }_\gamma }}\left[{\displaystyle \frac{3}{4}}\dot{\mathrm{\Delta }}_1+{\displaystyle \frac{\dot{a}}{a}}\mathrm{\Delta }_1\right].`$ (12)
Neglecting the expansion of the universe, equations (11a) and (12) can be combined to give a forced oscillator equation
$`\ddot{\mathrm{\Delta }}_0={\displaystyle \frac{k^2}{3R}}\mathrm{\Delta }_0+{\displaystyle \frac{2}{3}}\ddot{h},R1+{\displaystyle \frac{3\overline{\rho }_b}{4\overline{\rho }_\gamma }},`$ (13)
with solution
$`\mathrm{\Delta }_0(\tau )=\left(\mathrm{\Delta }_0(0){\displaystyle \frac{2R\ddot{h}}{k^2}}\right)\mathrm{cos}{\displaystyle \frac{k\tau }{\sqrt{3R}}}+{\displaystyle \frac{\sqrt{3R}}{k}}\dot{\mathrm{\Delta }}_0(0)\mathrm{sin}{\displaystyle \frac{k\tau }{\sqrt{3R}}}+{\displaystyle \frac{2R\ddot{h}}{k^2}},`$ (14)
where $`\mathrm{\Delta }_0(0)`$ and $`\dot{\mathrm{\Delta }}_0(0)`$ are evaluated when the wave first crosses the Hubble radius, $`k\tau 1`$. For adiabatic perturbations the first term dominates over the second because the perturbation breaks at $`k\tau 1`$ with $`\dot{\mathrm{\Delta }}_00`$. It is useful to define (gauge-invariant) radiation perturbation variables
$`\stackrel{~}{\mathrm{\Delta }}_0=\mathrm{\Delta }_0{\displaystyle \frac{2\ddot{h}}{k^2}},\stackrel{~}{\mathrm{\Delta }}_1=\mathrm{\Delta }_1+i{\displaystyle \frac{2\dot{h}}{3k}},`$
then the solution of equation (10) is
$`\stackrel{~}{\mathrm{\Delta }}(k,\mu ,\tau )={\displaystyle _0^\tau }\sigma _Tn_ea\left[\stackrel{~}{\mathrm{\Delta }}_0+4\mu \left(v_b+{\displaystyle \frac{i\dot{h}}{2k}}\right)\right]e^{ik\mu (\tau ^{}\tau )_\tau ^{}^\tau [\sigma _Tn_ea]𝑑\tau ^{\prime \prime }}𝑑\tau ^{}.`$ (15)
If $`k\tau 1`$, the second term in the square brackets is smaller than the first by a factor of $`k\tau `$, and the solution of equation (15) gives a power spectrum with a series of modulated acoustic peaks spaced at regular intervals of $`k_mr_s(a_r)=m\pi `$, where $`r_s`$ is the sound horizon at recombination
$`r_s={\displaystyle \frac{c}{\sqrt{3}H_0\mathrm{\Omega }_m^{1/2}}}{\displaystyle _0^{a_r}}{\displaystyle \frac{da}{\left(a+a_{equ}\right)^{1/2}}}{\displaystyle \frac{1}{R^{1/2}}},`$ (16)
(Hu and Sugiyama 1995). Here $`a_{equ}`$ is the scale factor when matter and radiation have equal densities and $`a_r`$ is the scale factor at recombination.
The multipole locations of the acoustic peaks in the angular power spectrum are given by
$`\mathrm{}_m=\alpha m\pi {\displaystyle \frac{d_A(z_r)}{r_s}}`$ (17)
where $`\alpha `$ is a number of order unity and $`d_A`$ is the angular diameter distance to last scattering
$`d_A={\displaystyle \frac{c}{H_0|\mathrm{\Omega }_k|^{1/2}}}\mathrm{sin}_k(|\mathrm{\Omega }_k|^{1/2}x)`$ (18a)
$`x{\displaystyle _{a_r}^1}{\displaystyle \frac{da}{[\mathrm{\Omega }_ma+\mathrm{\Omega }_ka^2+\mathrm{\Omega }_\mathrm{\Lambda }a^4]^{1/2}}}`$ (18b)
where $`\mathrm{\Omega }_k=1\mathrm{\Omega }_\mathrm{\Lambda }\mathrm{\Omega }_m`$ and $`\mathrm{sin}_k\mathrm{sinh}`$ if $`\mathrm{\Omega }_k>0`$ and $`\mathrm{sin}_k=\mathrm{sin}`$ if $`\mathrm{\Omega }_k<0`$.
The general dependence of the CMB power spectrum on cosmological parameters is therefore clear. The positions of the acoustic peaks depend on the geometry of the Universe via the angular diameter distance of equation (18) and on the value of the sound horizon $`r_s`$. The relative amplitudes of the peaks depend on the physical densities of the various constituents $`\omega _b\mathrm{\Omega }_bh^2`$, $`\omega _c\mathrm{\Omega }_ch^2`$, $`\omega _\nu \mathrm{\Omega }_\nu h^2`$, etc. and on the scalar fluctuation spectrum (parameterized here by a constant spectral index $`n_s`$). Clearly, models with the same initial fluctuation spectra and identical physical matter densities $`\omega _i`$ will have identical CMB power spectra at high multipoles if they have the same angular diameter distance to the last scattering surface. This leads to a strong geometrical degeneracy between $`\mathrm{\Omega }_m`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }`$ (e.g. Efstathiou and Bond 1999 and references therein). The power spectrum on large angular scales (equation 9) is sensitive to the spectral index and amplitude of the power spectrum, geometry of the Universe and, for extreme values of $`\mathrm{\Omega }_k`$ can break the geometrical degeneracy via the late-time Sachs-Wolfe effect. We will discuss briefly some of the constraints on cosmological parameters from the current CMB data in the next section. Before moving on to this topic, I mention some important points that cannot be covered in detail because of space limitations:
$``$ Inflationary models can give rise to tensor perturbations with a characteristic spectrum that declines sharply at $`\mathrm{}>100`$ (see e.g. Bond 1996 and references therein). In power-law like inflation, the tensor spectral index $`n_t`$ is closely linked to the scalar spectral index, $`n_tn_s1`$, and to the relative amplitude of the tensor and scalar perturbations.
$``$ The anisotropy of Thomson scattering causes the CMB anisotropies to be linearly polarized at the level of a few percent (see Bond 1996, Hu and White 1997, and references therein). Measurements of the linear polarization can distinguish between tensor and scalar perturbations and can constrain the epoch of reionization of the intergalactic medium (Zaldarriaga, Spergel and Seljak 1997).
$``$ The main effect of reionization is to depress the amplitude of the power spectrum at high multipoles by a factor of $`\mathrm{exp}(2\tau _{opt})`$ where $`\tau _{opt}`$ is the optical depth to Thomson scattering. In the ‘best fit’ CDM universe described in the next section ($`\omega _b=0.019`$, $`h=0.65`$, $`\mathrm{\Omega }_m=\mathrm{\Omega }_c+\mathrm{\Omega }_b0.3`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }0.7`$) and a reionization redshift of $`z_{reion}20`$ (a plausible value) $`\tau _{opt}0.2`$ which is significant. There is a reasonable chance that we might learn something about the ‘dark ages’ of cosmic history from precision measurements of the CMB.
## 4 Cosmological Parameters from the CMB
In this section, we review some of the constraints on cosmological parameters from the CMB data plotted in figure 1. The analysis is similar to that presented in Efstathiou et al. (1999, hereafter E99), in which we map the full likelihood function in $`5`$ parameters $`\mathrm{\Omega }_\mathrm{\Lambda }`$, $`\mathrm{\Omega }_m`$, $`\omega _c`$, $`n_s`$ and $`Q_{10}`$ (the amplitude of $`\sqrt{C_{\mathrm{}}}`$ at $`\mathrm{}=10`$ relative to that inferred from COBE). The baryon density is constrained to $`\omega _b=0.019`$, as determined from primordial nucleosynthesis and deuterium abundances measurements from quasar spectra (Burles and Tytler 1998). The results presented below are insensitive to modest variations ($`25\%`$) of $`\omega _b`$ and illustrate the main features of cosmological parameter estimation from the CMB. Recently, Tegmark and Zaldarriaga (2000) have performed a heroic $`10`$ parameter fit to the CMB data, including a tensor contribution and finite optical depth from reionization. I will discuss the effects of widening the parameter space briefly below but refer the reader to Tegmark and Zaldarriaga for a detailed analysis.
The best fit model in this five parameter space is plotted as the solid line in figure 1. It is encouraging that the best fitting model has perfectly reasonable parameters, a spatially flat universe with a nearly scale invariant fluctuation spectrum and a low CDM density $`\omega _c0.1`$. Marginalised likelihood functions are plotted in various projections in the parameter space in figures 2, 3 and 5 (uniform priors are assumed in computing the marginalized likelihoods, as described in E99). Figure 2 shows constraints on the position of the first acoustic peak measured by the ‘location’ parameter
$`\gamma _D={\displaystyle \frac{\mathrm{}_D(\mathrm{\Omega }_\mathrm{\Lambda },\mathrm{\Omega }_m)}{\mathrm{}_D(\mathrm{\Omega }_\mathrm{\Lambda }=0,\mathrm{\Omega }_m=1)}},`$ (19)
i.e. the parameter $`\gamma _D`$ measures the location of the acoustic peak relative to that in a spatially flat model with zero cosmological constant. The geometrical degeneracy between $`\mathrm{\Omega }_m`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }`$ described in the previous section is expressed by $`\gamma _D=\mathrm{constant}`$. Figure 2 shows that the best fitting value is $`\gamma _D=1`$ with a $`2\sigma `$ range of about $`\pm 0.3`$. The position of the first acoustic peak in the CMB data thus provides powerful evidence that the Universe is close to spatially flat.
Figure $`3`$ shows the marginalized likelihoods in the $`Q_{10}n_s`$ and $`\omega _cn_s`$ planes. The constraints on $`Q_{10}`$ and $`n_s`$ are not very different to those from the analysis of COBE alone (see e.g. Bond 1996). The experiments at higher multipoles are so degenerate with variations in other cosmological parameters that they do not help tighten the constraints on $`Q_{10}`$ and $`n_s`$. The constraints on $`\omega _c`$ and $`n_s`$ show an interesting result; if $`n_s1`$, then the best fit value of $`\omega _c`$ is about $`0.1`$ with a $`2\sigma `$ upper limit of about $`0.3`$. This constraint on $`\omega _c`$ comes from the height of the first acoustic peak, as shown in figure 4. In this diagram, the CMB data points have been averaged in $`10`$ band-power estimates as described by Bond, Knox and Jaffe (1998). The solid curve shows the best-fit model as plotted in figure 1, which has $`\omega _c=0.1`$. The dashed lines show models with $`\omega _c=0.25`$ and $`\omega _c=0.05`$ with the other parameters held fixed. Raising $`\omega _c`$ lowers the height of the peak and vice-versa. This result is not very sensitive to variations of $`\omega _b`$ in the neighbourhood of $`\omega _b0.02`$. Reionization and the addition of a tensor component can lower the height of the first peak relative to the anisotropies at lower multipoles and so the upper limits on $`\omega _c`$ are robust to the addition of these parameters. The CMB data have now reached the point where we have good constraints on the height of the first peak, as well as its location, and this is beginning to set interesting constraints on $`\omega _c`$. The best fit value of $`\mathrm{\Omega }_m0.3`$, derived from combining the CMB data with results from distant Type Ia supernovae (see figure 5) implies $`\omega _c0.11`$ for a Hubble constant of $`h=0.65`$, consistent with the low values of $`\omega _c`$ favoured by the height of the first acoustic peak.
The left hand panel of figure 5 shows the marginalized likelihood for the CMB data in the $`\mathrm{\Omega }_\mathrm{\Lambda }`$$`\mathrm{\Omega }_m`$ plane. The likelihood peaks along the line for spatially flat universes $`\mathrm{\Omega }_k=0`$ and it is interesting to compare with the equivalent figure in E99 to see how the new experimental results of the last year have caused the likelihood contours to narrow down around $`\mathrm{\Omega }_k=0`$. (See also Dodelson and Knox 1999 for a similar analysis using the latest CMB data). As is well known, the magnitude-redshift relation for distant Type Ia supernovae results in nearly orthogonal constraints in the $`\mathrm{\Omega }_\mathrm{\Lambda }`$$`\mathrm{\Omega }_m`$ plane, so combining the supernovae and CMB data can break the geometrical degeneracy. The right hand panel in Figure 5 combines the CMB likelihood function derived here with the likelihood function of the supernovae sample of Perlmutter et al. (1999) as analysed in E99. The combined likelihood function is peaked at $`\mathrm{\Omega }_m0.3`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }0.7`$.
It is remarkable how the CMB data and the supernovae data are homing in on a consistent set of cosmological parameters that are compatible with the simplest inflationary models and also with parameters inferred from a number of other observations (e.g. galaxy clustering, baryon content of clusters and dynamical estimates of the mean mass density, see Bahcall et al. 1999 for a review). It is also remarkable that the ‘best fit’ model requires a non-zero cosmological constant, a result that few cosmologists would have thought likely a few years ago.
The next few years will see a revolutionary increase in the volume and quality of CMB data. The results of the Boomerang Antarctic flight are awaited with great interest and should be of sufficient quality to render all previous analyses of cosmological parameters from the CMB obsolete. The polarization of the CMB has not yet been discovered, but a number of ground based and balloon borne experiments designed to detect polarization are under construction (Staggs, Gundersen and Church 1999). The MAP satellite, scheduled for launch in late 2000, will have polarization sensitivity and should determine the power spectrum $`C_{\mathrm{}}`$ accurately to about $`\mathrm{}800`$, defining the first three acoustic peaks. Further into the future, the Planck satellite, scheduled for launch in 2007, should determine the power spectrum to $`\mathrm{}>2500`$, provide sensitive polarization measurements and extremely accurate subtraction of foregrounds. Evidently, the era of precision cosmology is upon us and the next decade should see a dramatic improvement in our knowledge of fundamental cosmological parameters and in our understanding of the origin of fluctuations in the early Universe. |
warning/0002/hep-ph0002129.html | ar5iv | text | # Photon mixing in universes with large extra–dimensions
## I Introduction
It is well known that photons can be converted into gravitons by a magnetic field in a standard four dimensional spacetime. The propagation eigenstates are then mixtures of photon and graviton interaction eigenstates. In quantum words, this mixing is due to the fact that any particle which has a two–photon vertex can be created by a photon entering an external electromagnetic field and can oscillate coherently with the photon. Classically, this can be understood by the fact that an electromagnetic plane wave cannot radiate gravitationally in vacuum since its stress–energy tensor contains no quadrupole . But, a time varying quadrupole appears (due to interference) when an electromagnetic plane wave propagates through a constant magnetic field . The implications of this effect on the cosmic microwave background (photons) has been considered and it has been shown that it will be undetectable for standard cosmological magnetic fields . A similar effect also generically happens for axions (and for any particle having a two–photon vertex). The photon–axion (see for reviews on axions) mixing has yet been studied in details by many authors (see e.g. ) and is used in experiments, since the pioneer work by Sikivie , to put constraints on the axion parameters (see e.g. for an up to date review on such experiments).
Recently, a lot of interest has been raised by models where the universe has large extra–dimensions . In such models, the Planck scale, $`M_4`$, is no longer a fundamental scale but is related to the fundamental mass scale of the $`D`$ dimensional theory, $`M_D`$, through
$$\overline{M}_4^2R^nM_D^{n+2},$$
(1)
where $`R`$ is a length scale (usually taken to be the radius of the $`n=D4`$ compact extra–dimensions) and $`\overline{M}_4M_4/\sqrt{8\pi }`$. $`M_D`$ can be significantly smaller than $`\overline{M}_4`$ at the price of having large extra–dimensions. These ideas can be naturally embedded in fundamental string theories with a low string scale (see also for earlier discussions on TeV scale extra–dimensions). In these models, gravity can propagate in the $`D`$ dimensional spacetime (bulk space time) whereas the standard model fields are localised on a 3–brane. An effect of the compact extra–dimensions arises from interactions between the Kaluza–Klein (KK) excitations of the gravitons (or other bulk fields) which are seen in four dimensions as a tower of massive particles . Constraints on the size of these extra–dimensions can be obtained both from the laboratory physics and from astrophysics and cosmology (see e.g. ). For instance, the emission of KK gravitons induces an energy loss in many astrophysical objects such as the Sun, red giants and supernovae SN1987A implying the lower bound $`M_D>30130\text{Tev}`$ (2.1 - 9.2) TeV for the case of $`n=2`$ (3) extra–dimensions . Some authors also have recently investigated the possible presence of axions in the bulk which would be coupled to the brane degrees of freedom. Such a bulk axion also gives rise to a tower of KK-states as seen from a four dimensional point of view.
The goal of this article is to investigate the effects of the photon–KK graviton and photon–KK axion oscillations and to estimate their effects in cosmology and astrophysics, as well as terrestrial experiments. Since there is a large number of KK states with which the photon can mix, one can expect a departure from the usual four dimensional result. We first describe (§ II) the photon–KK graviton system starting from a $`D`$ dimensional action, linearising it and compactifying it to four dimensions. We also show (§ III) that the photon–axion mixing is described by the same formalism and lead to the same effects as for gravitons. Then, we turn to investigate the mixing in itself. For that purpose, we sum up the known results of the mixing of a photon with a low mass particle in four dimensions (§ IV A) and then discuss the most general case of a $`D`$ dimensional spacetime (§ IV B). The general expression for the oscillation probability is then evaluated in the particular cases of a five (§ V) and of a six (§ VI) dimensional spacetime. We show that, as long as one is in a weak coupling regime, one can add the individual probabilities which leads to an enhancement of the oscillation probability if the KK modes are light enough. We also show that there exists a regime where the photon mixes strongly preferentially with a given KK mode. We generalise our results to an inhomogeneous magnetic field (§ VII) when the probability of oscillation is small, and then turn to the cosmological and astrophysical situations where such effects may be observed. We study the case of the cosmic microwave background (§ VIII A), of pulsars (§ VIII B) and of magnetars (§ VIII C). We show that even if the enhancement of the oscillation probability can be very important, it is still very difficult to observe this effect in known astrophysical systems. To finish (§ IX) we discuss laboratory experiments and particularly polarisation experiments. For that purpose, we describe the computation of the phase shift between the two polarisations of an electromagnetic wave and compare the result to the standard four dimensional case.
## II Equations of motion for the photon–graviton system
Following , we consider a field theory defined by the $`D`$ dimensional action
$$S_D=\frac{1}{2\kappa _D^2}d^Dz\sqrt{\overline{g}}\overline{R}+d^Dz\sqrt{\overline{g}}_m,$$
(2)
where $`\overline{g}_{AB}`$ is the $`D`$ dimensional metric with signature $`(,+,\mathrm{},+)`$, $`\kappa _D^28\pi G_D=\overline{M}_D^{(2+n)}`$ and $`_m`$ is the matter Lagrangian. The indices $`A,B,\mathrm{}`$ take the value $`0,\mathrm{..3},5..,D`$ and we decompose $`z^A`$ as
$$z^A=(x^\mu ,y^a)$$
(3)
with $`\mu ,\nu ,\mathrm{}=0,\mathrm{},3`$ and $`a,b,\mathrm{}=5,\mathrm{},D`$.
This theory is considered as being a low energy effective theory valid below some cut–off $`M_{\mathrm{max}}`$ in energy (see e.g. ). We will discuss in this paper only the cases $`n=1`$ and $`n=2`$ in details, and our conclusions, in these cases, are mostly cut–off independent. We stress that the relationship between this cut–off and the fundamental string scale (if one wishes to embed these theories in superstring models) can be much more complicated than what is naively expected (see ).
We expand the metric around the $`D`$ dimensional Minkowski spacetime as
$$g_{AB}=\eta _{AB}+\frac{h_{AB}}{\overline{M}_D^{1+n/2}}$$
(4)
where $`\eta _{AB}`$ is the $`D`$ dimensional Minkowski metric. Inserting (4) in (2) and using the definition of the stress–energy tensor as $`\sqrt{\overline{g}}T_{AB}2\delta (_m\sqrt{\overline{g}})/\delta g^{AB}`$, so that $`_m\sqrt{\overline{g}}=_{m0}h^{AB}T_{AB}/2\overline{M}_D^{1+n/2}`$, we obtain the linearised action
$`S_D`$ $`=`$ $`{\displaystyle }d^Dz[{\displaystyle \frac{1}{2}}h^{AB}^C_Ch_{AB}{\displaystyle \frac{1}{2}}h_A^A^C_Ch_B^B+{\displaystyle \frac{1}{2}}h^{AB}_A_Bh_C^C+{\displaystyle \frac{1}{2}}h_A^A_C_Bh^{CB}`$ (6)
$`h^{AB}_A_Ch_B^C{\displaystyle \frac{1}{\overline{M}_D^{1+n/2}}}h^{AB}T_{AB}+_{m0}].`$
We compactify this theory to get a four dimensional theory and use the periodicity on $`y^a`$ to expand the field $`h^{AB}`$ as
$$h_{AB}(z^A)=\underset{\stackrel{}{p}Z}{}\frac{h_{AB}^{(\stackrel{}{p})}(x^\mu )}{\sqrt{V_n}}\mathrm{exp}(i\frac{p^ay_a}{R})$$
(7)
where $`V_n=(2\pi R)^n`$ is the volume of the compact $`n`$ dimensional space (assumed to be a cubic $`n`$–torus). $`h_{AB}`$ is split into a sum of KK modes living in the four dimensional spacetime. The ordinary matter being confined to the brane, its stress–energy tensor must satisfy
$$T_{AB}(z^C)=T_{\mu \nu }(x^\lambda )\delta ^{(n)}(y^c)\eta _A^\mu \eta _B^\nu .$$
(8)
In which follows, we restrict our attention to the case of an electromagnetic field $`F_{\mu \nu }`$ for which the stress–energy tensor is given by
$$T_{\mu \nu }=F_{\mu \lambda }F_\nu ^\lambda \frac{1}{4}\eta _{\mu \nu }F^{\lambda \rho }F_{\lambda \rho }.$$
(9)
The fields $`h_{AB}^{(\stackrel{}{p})}`$ can be decomposed into spin–2, spin–1 and spin–0 four dimensional fields . Only spin–2 and spin–0 particles couple to ordinary matter and spin–0 particles couple only to $`T_\lambda ^\lambda `$. For the electromagnetic field $`T_\lambda ^\lambda =0`$ classically so that the only relevant KK modes (at the tree level analysis of this article) will be the spin–2 particles $`G_{\mu \nu }^{(\stackrel{}{p})}`$ for which the action (6) reduces to
$`S_4`$ $`=`$ $`{\displaystyle }d^4x[{\displaystyle \frac{1}{2}}G_{\mu \nu }^{(\stackrel{}{p})}(\mathrm{}m_\stackrel{}{p}^2)G^{(\stackrel{}{p})\mu \nu }+G^{(\stackrel{}{p})\mu \nu }_\mu _\lambda G_\nu ^{(\stackrel{}{p})\lambda }{\displaystyle \frac{1}{2}}G_\mu ^{(\stackrel{}{p})\mu }(\mathrm{}m_\stackrel{}{p}^2)G_\nu ^{(\stackrel{}{p})\nu }G^{(\stackrel{}{p})\mu \nu }_\mu _\nu G_\lambda ^{(\stackrel{}{p})\lambda }`$ (11)
$`{\displaystyle \frac{1}{\overline{M}_4}}G^{(\stackrel{}{p})\mu \nu }T_{\mu \nu }{\displaystyle \frac{1}{4}}F^{\mu \nu }F_{\mu \nu }],`$
where $`m_\stackrel{}{p}^2=\stackrel{}{p}^2/R^2`$ is the square mass of the KK graviton, $`\mathrm{}_\mu ^\mu `$.
The equations of motion deduced from the Lagrangian (11) are the coupled Einstein–Maxwell equations
$`(\mathrm{}m_\stackrel{}{p}^2)G_{\mu \nu }^{(\stackrel{}{p})}={\displaystyle \frac{2}{\overline{M}_4}}T_{\mu \nu }`$ (12)
$`^\mu G_{\mu \nu }^{(\stackrel{}{p})}=0`$ (13)
$`G_\mu ^{(\stackrel{}{p})\mu }=0`$ (14)
$`_\alpha F^{\alpha \beta }{\displaystyle \frac{2}{\overline{M}_4}}{\displaystyle \underset{\stackrel{}{p}}{}}_\alpha \left(G^{(\stackrel{}{p})\alpha \nu }F_\nu ^\beta G^{(\stackrel{}{p})\beta \nu }F_\nu ^\alpha \right)=0.`$ (15)
When $`\stackrel{}{p}0`$, the field $`G_{\mu \nu }^{(\stackrel{}{p})}`$ has 10-1-4=5 degrees of freedom which is what is expected for a massive spin–2 particle.
We now consider an electromagnetic plane wave in the presence of a magnetic field $`\stackrel{}{H}_0`$ which is assumed constant on a characteristic scale $`\mathrm{\Lambda }_c`$ in the sense that its variation in space and time are negligible on scales comparable to the photon wavelength and period. We define the basis
$$\stackrel{}{e}_{}\frac{\stackrel{}{k}}{k},\stackrel{}{e}_\times \frac{\stackrel{}{H}_0}{H_0},\stackrel{}{e}_+,$$
(16)
such that $`(\stackrel{}{e}_\times ,\stackrel{}{e}_+,\stackrel{}{e}_{})`$ is a direct orthonormal basis of the three dimensional space and where $`\stackrel{}{H}_0`$ is the perpendicular component of $`\stackrel{}{H}_0`$ with respect to the direction of propagation $`\stackrel{}{k}`$. We decompose the KK gravitons in scalar (S), vector (V) and tensor (T) as
$`(S)`$ $`G_{00}^{(\stackrel{}{p})}=\varphi ^{(\stackrel{}{p})},G_{0i}^{(\stackrel{}{p})}=ik_ik^2\dot{\varphi }^{(\stackrel{}{p})},G_{ij}^{(\stackrel{}{p})}={\displaystyle \frac{\varphi ^{(\stackrel{}{p})}}{3}}\delta _{ij}{\displaystyle \frac{3}{2k^2}}\mathrm{\Delta }_{ij}\left({\displaystyle \frac{\varphi }{3}}+{\displaystyle \frac{\ddot{\varphi }}{k^2}}\right)^{(\stackrel{}{p})}`$ (17)
$`(V)`$ $`G_{0i}^{(\stackrel{}{p})}V_i^{(\stackrel{}{p})}=V_+^{(\stackrel{}{p})}e_i^++V_\times ^{(\stackrel{}{p})}e_i^\times ,G_{00}^{(\stackrel{}{p})}=0,G_{ij}^{(\stackrel{}{p})}={\displaystyle \frac{2}{k^2}}k_{(j}\dot{V}_{i)}^{(\stackrel{}{p})}`$ (18)
$`(T)`$ $`G_{00}^{(\stackrel{}{p})}=0,G_{0i}^{(\stackrel{}{p})}=0,G_{ij}^{(\stackrel{}{p})}=G_+^{(\stackrel{}{p})}ϵ_{ij}^++G_\times ^{(\stackrel{}{p})}ϵ_{ij}^\times ,`$ (19)
where $`\mathrm{\Delta }_{ij}\left(k_ik_j\frac{k^2}{3}\delta _{ij}\right)`$, $`i,j=\mathrm{1..3}`$ and a dot refers to a time derivative. The polarisation tensor of the graviton modes $`ϵ_{ij}^\lambda `$ is defined by
$`ϵ_{ij}^\lambda \left(e_i^\times e_j^\times e_i^+e_j^+\right)\delta _\times ^\lambda +2e_i^{(+}e_j^{\times )}\delta _+^\lambda .`$ (20)
The advantage of such a decomposition is that the scalar, vector and tensor contributions decouple. The five degrees of freedom of each massive spin–2 KK gravitons have been decomposed in one scalar mode ($`\varphi ^{(\stackrel{}{p})}`$), two vector modes ($`V_{+/\times }^{(\stackrel{}{p})}`$) and two tensor modes ($`G_{+/\times }^{(\stackrel{}{p})}`$). Each of these modes satisfies independently the constraints (1314).
We consider an electromagnetic wave with a potential vector of the form
$$\stackrel{}{A}=i(A_\times (u),A_+(u),0)\text{e}^{i\omega t},$$
(21)
where $`u`$ is the coordinate along the direction of propagation. We have introduced the arbitrary phase $`i`$ so that its electric and magnetic fields are
$`\stackrel{}{E}`$ $``$ $`_t\stackrel{}{A}=(\omega A_\times (u),\omega A_+(u),0)\text{e}^{i\omega t}`$ (22)
$`\stackrel{}{B}`$ $``$ $`\text{curl}(\stackrel{}{A})=(i_uA_+(u),i_uA_\times (u),0)\text{e}^{i\omega t}.`$ (23)
The stress–energy tensor of these waves in the presence of $`\stackrel{}{H}_0`$ has no vector component. Its tensor component is given by
$$T_{ij}=i\underset{\lambda =+,\times }{}_uA_\lambda H_0\text{e}^{i\omega t}ϵ_{ij}^\lambda .$$
(24)
We see, as expected, that a plane wave possesses a tensor part only if it propagates in an external field and that the polarisations $`+`$ and $`\times `$ of the electromagnetic wave couple respectively to the polarisations $`+`$ and $`\times `$ of the gravitons. The electromagnetic wave generates also scalar perturbations, but it can be shown that (in the usual four dimensional case) the total energy converted in this scalar wave are negligible compared to the tensor contribution. In the following, we concentrate on the tensor modes.
The equation of evolution of this system is given by the Einstein equation (12) which reduces to
$$\left(\omega ^2+_u^2m_\stackrel{}{p}^2\right)G_\lambda ^{(\stackrel{}{p})}=\frac{2iH_0}{\overline{M}_4}_uA_\lambda ,$$
(25)
and the Maxwell equation (15) which reduces to
$$\left(\omega ^2+_u^2\right)A_\lambda =\frac{2iH_0}{\overline{M}_4}\underset{\stackrel{}{p}}{}_uG_\lambda ^{(\stackrel{}{p})},$$
(26)
where we have used the ansatz $`G_{ij}=_\lambda G_\lambda (u)\text{e}^{i\omega t}ϵ_{ij}^\lambda `$ for the gravitons.
Since we have assumed that the magnetic field varies in space on scales much larger than the photon wavelength, we can perform the expansion $`\omega ^2+_u^2=(\omega +i_u)(\omega i_u)=(\omega +k)(\omega i_u)`$ for a field propagating in the $`+u`$ direction. If we assume a general dispersion equation of the form $`\omega =nk`$ and that the refractive index $`n`$ satisfies $`|n1|1`$, we may approximate $`\omega +k=2\omega `$ and $`k/\omega =1`$. This approximation can be understood as a WKB limit where we set $`A(u)=|A(u)|\text{e}^{iku}`$ and assume that the amplitude $`|A|`$ varies slowly, i.e. that $`_u|A|k|A|`$. In that limit, the system (2526) reduces to
$$\left[\omega i_u+_\lambda \right]\left[\begin{array}{c}A_\lambda \\ G_\lambda ^{(0)}\\ \mathrm{}\\ G_\lambda ^{(q)}\\ \mathrm{}\end{array}\right]=0,$$
(27)
the matrix $`_\lambda `$ being given by
$$\left(\begin{array}{ccccccc}\mathrm{\Delta }_\lambda & \mathrm{\Delta }_M& \mathrm{\Delta }_M& \mathrm{}& \mathrm{\Delta }_M& \mathrm{}& \mathrm{}\\ \mathrm{\Delta }_M& \mathrm{\Delta }_m^{(0)}& 0& \mathrm{}& 0& \mathrm{}& \mathrm{}\\ \mathrm{\Delta }_M& 0& \mathrm{\Delta }_m^{(1)}& 0& 0& \mathrm{}& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ \mathrm{\Delta }_M& 0& \mathrm{}& 0& \mathrm{\Delta }_m^{(q)}& 0& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\end{array}\right)$$
(28)
with
$$\mathrm{\Delta }_M\frac{H_0}{\overline{M}_4}\text{and}\mathrm{\Delta }_m^{(q)}\stackrel{}{p}_{(q)}^2\mathrm{\Delta }_m.$$
(29)
$`\stackrel{}{p}_{(q)}`$ is a n–uplets $`(p_{q_1},p_{q_2},..p_{q_n})`$ of integers and we have ordered the $`\mathrm{\Delta }_m^{(q)}`$ such that
$$\left|\mathrm{\Delta }_m^{(q)}\right|\left|\mathrm{\Delta }_m^{(q+1)}\right|$$
and $`\mathrm{\Delta }_m`$ is defined by
$$\mathrm{\Delta }_m\frac{1}{2R^2\omega }.$$
(30)
Each $`\mathrm{\Delta }_m^{(q)}`$ appears with a multiplicity given by the number of n–uplets having the same norm $`_{i=1..n}p_{q_i}^2`$. We define the two series $`(r_i)_{i1}`$ and $`(s_i)_{i1}`$ such that
$$\mathrm{\Delta }_m^{(r_i1)}<\mathrm{\Delta }_m^{(r_i)}=\mathrm{\Delta }_m^{(r_i+1)}=\mathrm{}\mathrm{\Delta }_m^{(r_i+s_i1)}<\mathrm{\Delta }_m^{(r_i+s_i)}\mathrm{\Delta }_m^{(r_{i+1})}.$$
(31)
We have $`r_{i+1}=r_i+s_i`$, and $`s_i`$ is the multiplicity of the element $`\mathrm{\Delta }_m^{(r_i)}`$, i.e. the number of times it appears in the matrix (28). $`r_i`$ is the rank in the series $`(\mathrm{\Delta }_m^{(0)},\mathrm{\Delta }_m^{(1)},\mathrm{})`$ where the $`i^{\mathrm{th}}`$ distinct value of $`\mathrm{\Delta }_m^{(q)}`$ appears for the first time. In the case of a five dimensional spacetime one can easily find out that $`s_1=1`$ and $`s_i=2`$ for $`i>1`$ and that $`r_1=0,r_2=1,r_3=3,\mathrm{}`$ In the case of a six dimensional spacetime, $`s_i=(1,4,4,\mathrm{})`$, $`r_i=(0,1,5,9,\mathrm{})`$. Introducing the cut–off $`M_{\mathrm{max}}`$ discussed above, we require $`m_\stackrel{}{p}^2=\stackrel{}{p}^2/R^2<M_{\mathrm{max}}^2`$, which using (1) translates into $`\stackrel{}{p}^2<p_{\mathrm{max}}^2`$ with
$$p_{\mathrm{max}}=\left(\frac{\overline{M}_4}{M_D}\right)^{2/n}\left(\frac{M_{\mathrm{max}}}{M_D}\right).$$
(32)
Setting $`M_{\mathrm{max}}M_D`$, we obtain $`p_{\mathrm{max}}\left(\overline{M}_4/M_D\right)^{2/n}`$. For $`n=2`$ and $`M_D1\mathrm{T}\mathrm{e}\mathrm{V}`$ one finds
$$p_{\mathrm{max}}10^{15},$$
(33)
which means that we have to consider a very large number of KK states. We also define a maximum index, $`N`$ say, for the series $`\mathrm{\Delta }_m^{(q)}`$ defined by
$$N\text{sup}\{q|\stackrel{}{p}_{(q)}^2=p_{\mathrm{max}}^2\},$$
(34)
which translates into a maximum index $`N_D`$ for the series $`s_i`$ and $`r_i`$. We stress that the number of KK modes relevant for the photon– KK graviton oscillation is likely to be smaller than $`N`$ due to decoherence effects, such as the source and detector finite width in momentum, the wave packet separation for massive (and non–relativistic) KK modes… (see e.g. for a description of these effects in the case of neutrino oscillation).
The term $`\mathrm{\Delta }_\lambda `$ can be decomposed as $`\mathrm{\Delta }_\lambda =\mathrm{\Delta }_{\mathrm{QED}}+\mathrm{\Delta }_{\mathrm{CM}}+\mathrm{\Delta }_{\mathrm{plasma}}`$. The first term contains the effect of vacuum polarisation giving a refractive index to the photon (see e.g. Adler ) and can be computed by adding the Euler–Heisenberg effective Lagrangian which is the lowest order term of the non–linearity of the Maxwell equations in vacuum (see e.g. ) to the action (11) The equation of motion derived from (11) is (27) with $`\mathrm{\Delta }_\lambda =0`$. We intentionaly omit the Euler–Heisenberg contribution in the presentation for the sake of clarity. Its Lagrangian is explicitely given by $`_{EH}=\frac{\alpha ^2}{90m_e^4}\left[(F^{\mu \nu }F_{\mu \nu })^2+\frac{7}{4}(F^{\mu \nu }\stackrel{~}{F}_{\mu \nu })^2\right]`$.. The second term describes the Cotton–Mouton effect, i.e. the birefringence of gases and liquids in presence of a magnetic field and the third term the effect of the plasma (since, in general, the photon does not propagate in vacuum). Their explicit expressions are given by
$`\mathrm{\Delta }_{\mathrm{QED}}^\times ={\displaystyle \frac{7}{2}}\omega \xi ,\mathrm{\Delta }_{\mathrm{QED}}^+=2\omega \xi ,`$ (35)
$`\mathrm{\Delta }_{\mathrm{plasma}}={\displaystyle \frac{\omega _{\mathrm{plasma}}^2}{2\omega }},`$ (36)
$`\mathrm{\Delta }_{\mathrm{CM}}^\times \mathrm{\Delta }_{\mathrm{CM}}^+=2\pi CH_0^2,`$ (37)
with $`\xi (\alpha /45\pi )(H_0/H_c)^2`$, $`H_cm_e^2/e=4.41\times 10^{13}\text{G}`$, $`m_e`$ the electron mass, $`e`$ the electron charge and $`\alpha `$ the fine structure constant. $`C`$ is the Cotton–Mouton constant ; this effect gives only the difference of the refractive indices and the exact value of $`C`$ is hard to determine ; we will neglect this effect but for the polarisation experiments (see § IX B). The plasma frequency $`\omega _{\mathrm{plasma}}`$ is defined by
$$\omega _{\mathrm{plasma}}^24\pi \alpha \frac{n_e}{m_e},$$
(38)
$`n_e`$ being the electron density. Note that the $`\mathrm{\Delta }_m^{(q)}`$ are always negative whereas $`\mathrm{\Delta }_\lambda `$ is positive if the contribution of the vacuum dominates and negative when the plasma term dominates.
The equation of motion (27) reduces to the one studied by Raffelt and Stodolsky when one considers only four dimensions so that $`_\lambda `$ contains only the massless graviton and is then a $`2\times 2`$ matrix. The main difference lies in the fact that now the electromagnetic component couples to a large numbers of KK gravitons. This can be compared to some models of neutrino oscillations in spacetime with extra–dimensions . We should also note that the two polarisations are, as expected, completely decoupled and obey the same equation of evolution. In the following of this article, we omit the subscript $`\lambda `$ of the polarisation.
## III Equations of motion for the photon–axion system
Before turning to the study the photon–graviton mixing, we consider the case of axions and show that the photon–axion mixing can be described by the same formalism.
We consider the generic action for the bulk axion photon system
$`S_4`$ $`=`$ $`{\displaystyle d^4x\left[\underset{\stackrel{}{p}}{}\left(\frac{1}{2}\left\{^\mu a^{(\stackrel{}{p})}_\mu a^{(\stackrel{}{p})}+m_\stackrel{}{p}^2a_{}^{(\stackrel{}{p})}{}_{}{}^{2}\right\}+\frac{a^{(\stackrel{}{p})}}{f_{\mathrm{PQ}}}F_{\mu \nu }\stackrel{~}{F}^{\mu \nu }\right)\frac{1}{4}F_{\mu \nu }F^{\mu \nu }\right]},`$ (39)
where the $`a^{(\stackrel{}{p})}`$ are the mass eigenstates of the axions and $`m_\stackrel{}{p}`$ their masses. $`\stackrel{~}{F}_{\mu \nu }\frac{1}{2}ϵ_{\mu \nu \rho \sigma }F^{\rho \sigma }`$ is the dual of the electromagnetic tensor, $`ϵ_{\mu \nu \rho \sigma }`$ being the completely antisymetric tensor such that $`ϵ_{0123}=+1`$. As for the gravitons, the mass spectrum is expected to be discrete, the states can be labelled by a n-uplet $`\stackrel{}{p}`$ and is expected to have a typical spacing of $`1/R`$. We have considered here that every axion KK state $`a^{(\stackrel{}{p})}`$ couples to the photon with the same coupling $`1/f_{\mathrm{PQ}}`$. This is only expected to be true if the typical mass, $`m_{\mathrm{PQ}}`$, given to the axion zero mode by instanton effects is much lower than the typical KK mass $`1/R`$ . Let us further stress here that for such bulk axions the usual relationship between the axion mass and the PQ scale does not hold anymore, so that one expects to see interesting new effects to appear . Inspired by the usual bounds on $`f_{\mathrm{PQ}}`$, we take $`f_{\mathrm{PQ}}`$ of order $`10^{10}`$ GeV. However we stress that the usual bounds on $`f_{\mathrm{PQ}}`$ may be modified partly because of a large number of axion–like particles coupling to the standard model fields. For example, we expect that the astrophysical bounds will be more stringent mainly because a star will now be able to emit all the energetically accessible modes (see and also for a discussion on relic axions oscillations).
We do not consider the perturbations of the metric and work in Minkowski spacetime since we are interested in the interaction between the photon and the axion. We deduce from (39) the coupled Klein–Gordon and Maxwell equations
$`\left(\mathrm{}m_\stackrel{}{p}^2\right)a^{(\stackrel{}{p})}={\displaystyle \frac{1}{f_{\mathrm{PQ}}}}F_{\mu \nu }\stackrel{~}{F}^{\mu \nu },`$ (40)
$`_\alpha F^{\alpha \beta }={\displaystyle \frac{4}{f_{\mathrm{PQ}}}}_\alpha \left[{\displaystyle \underset{\stackrel{}{p}}{}}a^{(\stackrel{}{p})}\stackrel{~}{F}^{\alpha \beta }\right].`$ (41)
We now decompose the electromagnetic wave as in (21-23) with respect to the basis (16), so that the former system reads
$`\left(\mathrm{}m_\stackrel{}{p}^2\right)a^{(\stackrel{}{p})}={\displaystyle \frac{4H_0}{f_{\mathrm{PQ}}}}A_\times `$ (42)
$`\mathrm{}A_\lambda ={\displaystyle \frac{4H_0}{f_{\mathrm{PQ}}}}\omega \delta _{\lambda \times }{\displaystyle \underset{\stackrel{}{p}}{}}a^{(\stackrel{}{p})},`$ (43)
where we have decomposed the axions as $`a^{(\stackrel{}{p})}(u)\mathrm{exp}(i\omega t)`$. Using the same WKB limit as in section IV, we obtain the linearised system
$$\left(\omega i_u+\right)\left[\begin{array}{ccc}A_+& & \\ A_\times & & \\ a^{(\stackrel{}{0})}& & \\ \mathrm{}& & \\ a^{(\stackrel{}{p})}& & \\ \mathrm{}& & \end{array}\right]=0.$$
(44)
The matrix $``$ is now defined by
$$=\left(\begin{array}{cccccc}\mathrm{\Delta }_+& 0& 0& 0& 0& \mathrm{}\\ 0& \mathrm{\Delta }_\times & \mathrm{\Delta }_M& \mathrm{\Delta }_M& \mathrm{}& \\ 0& \mathrm{\Delta }_M& \mathrm{\Delta }_a^{(0)}& 0& 0& \mathrm{}\\ 0& \mathrm{\Delta }_M& 0& \mathrm{}& \mathrm{}& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{\Delta }_a^{(q)}& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\end{array}\right)$$
(45)
with
$$\mathrm{\Delta }_M=\frac{4H_0}{f_{\mathrm{PQ}}},\mathrm{\Delta }_a^{(\stackrel{}{p})}=\frac{m_\stackrel{}{p}^2}{2\omega },$$
(46)
$`\mathrm{\Delta }_+`$ and $`\mathrm{\Delta }_\times `$ being given by equation (35). This system reduces to the Raffelt and Stodolsky system when we consider only four dimensions. Only the component $`\times `$, i.e. parallel to the magnetic field, couples to the axions. This is a major difference compared with gravitons for which both polarisations of the photon evolve alike whereas here only $`A_\times `$ is affected by the mixing.
One of the goal of this section was to set the theoretical framework for further experimental studies of photon–bulk axion oscillations (see § IX for a more detailed discussion) and to show that it is described by a similar formalism as photon–KK graviton oscillations (under the validity conditions explained below equation (39).
## IV Photon–KK state mixing in a homogeneous field
We now describe the physical implications of the system (2526) and start by reviewing briefly the well studied problem of the mixing of a photon with a low mass particle in four dimensions (§ IV A). We then give the exact expression of the oscillation probability in $`D`$ dimensions and discuss qualitatively its magnitude and the effect of the coupling of the photon to a large number of particles.
### A The usual photon mixing with a low mass particle
This case was well studied in the literature, see e.g. Raffelt and Stodolsky and we just summarize the main features of the results to compare to the case of a spacetime with extra–dimensions. For the mixing with a single particle of mass $`m`$, the matrix $``$ reduces to
$$=\left(\begin{array}{cc}\mathrm{\Delta }_\lambda & \mathrm{\Delta }_M\\ \mathrm{\Delta }_M& \mathrm{\Delta }_m\end{array}\right)$$
(47)
with $`\mathrm{\Delta }_mm^2/2\omega `$. The solution to the equation of motion (27) is obtained by diagonalising $``$ throught a rotation
$$\left[\begin{array}{c}A^{}\\ G^{}\end{array}\right]=\left(\begin{array}{cc}\mathrm{cos}\vartheta & \mathrm{sin}\vartheta \\ \mathrm{sin}\vartheta & \mathrm{cos}\vartheta \end{array}\right)\left[\begin{array}{c}A\\ G\end{array}\right]$$
(48)
with
$$\mathrm{tan}2\vartheta 2\frac{\mathrm{\Delta }_M}{\mathrm{\Delta }_\lambda \mathrm{\Delta }_m}=\frac{2\alpha }{1\beta }$$
(49)
and where we have introduced $`\alpha \mathrm{\Delta }_M/\mathrm{\Delta }_\lambda `$ and $`\beta \mathrm{\Delta }_m/\mathrm{\Delta }_\lambda `$. We obtain by solving (27) in this new basis
$`A^{}(u)`$ $`=`$ $`\text{e}^{i\mathrm{\Delta }_\lambda ^{}u}A^{}(0)`$ (50)
$`G^{}(u)`$ $`=`$ $`\text{e}^{i\mathrm{\Delta }_g^{}u}G^{}(0)`$ (51)
where a global phase $`\omega u`$ has been omitted. The two eigenvalues $`\mathrm{\Delta }_\lambda ^{}`$ and $`\mathrm{\Delta }_g^{}`$ of $``$ are explicitly given by
$`\mathrm{\Delta }_\lambda ^{}={\displaystyle \frac{\mathrm{\Delta }_\lambda +\mathrm{\Delta }_m}{2}}+{\displaystyle \frac{\mathrm{\Delta }_\lambda \mathrm{\Delta }_m}{2\mathrm{cos}2\vartheta }}\mathrm{and}\mathrm{\Delta }_g^{}={\displaystyle \frac{\mathrm{\Delta }_\lambda +\mathrm{\Delta }_m}{2}}{\displaystyle \frac{\mathrm{\Delta }_\lambda \mathrm{\Delta }_m}{2\mathrm{cos}2\vartheta }}.`$ (52)
Going back to the initial basis, we obtain
$`A(u)`$ $`=`$ $`\left(\text{e}^{i\mathrm{\Delta }_\lambda ^{}u}\mathrm{cos}^2\vartheta +\text{e}^{i\mathrm{\Delta }_g^{}u}\mathrm{sin}^2\vartheta \right)A(0)+\mathrm{sin}\vartheta \mathrm{cos}\vartheta \left(\text{e}^{i\mathrm{\Delta }_\lambda ^{}u}\text{e}^{i\mathrm{\Delta }_g^{}u}\right)G(0),`$ (53)
$`G(u)`$ $`=`$ $`\mathrm{sin}\vartheta \mathrm{cos}\vartheta \left(\text{e}^{i\mathrm{\Delta }_\lambda ^{}u}\text{e}^{i\mathrm{\Delta }_g^{}u}\right)A(0)+\left(\text{e}^{i\mathrm{\Delta }_g^{}u}\mathrm{cos}^2\vartheta +\text{e}^{i\mathrm{\Delta }_\lambda ^{}u}\mathrm{sin}^2\vartheta \right)G(0).`$ (54)
The oscillation probability of a photon into a graviton is computed by considering the initial state $`(A(0)=1,G(0)=0)`$ and is given by
$`P(\gamma g)A(0)G(u)^2`$ $`=`$ $`\mathrm{sin}^2\left(2\vartheta \right)\mathrm{sin}^2\left({\displaystyle \frac{\mathrm{\Delta }_{\mathrm{osc}}}{2}}u\right),`$ (55)
$`=`$ $`\left(\mathrm{\Delta }_Mu\right)^2{\displaystyle \frac{\mathrm{sin}^2(\mathrm{\Delta }_{\mathrm{osc}}u/2)}{(\mathrm{\Delta }_{\mathrm{osc}}u/2)^2}}`$ (56)
with
$$\mathrm{\Delta }_{\mathrm{osc}}\mathrm{\Delta }_\lambda ^{}\mathrm{\Delta }_g^{}=\frac{\mathrm{\Delta }_\lambda \mathrm{\Delta }_m}{\mathrm{cos}2\vartheta }=\frac{2\mathrm{\Delta }_M}{\mathrm{sin}2\vartheta }=\frac{1\beta }{\mathrm{cos}2\vartheta }\mathrm{\Delta }_\lambda $$
(57)
The oscillation length is thus given by $`\mathrm{}_{\mathrm{osc}}2\pi /\mathrm{\Delta }_{\mathrm{osc}}`$. We see that the oscillation probability cannot be larger than $`\left(\mathrm{\Delta }_Mu\right)^2`$. One has to be aware that $`\vartheta `$ depends on the polarisation index $`\lambda `$.
It is interesting to single out the two following limiting regimes:
* The weak mixing regime in which $`\vartheta 1`$ so that the probability (55) reduces to
$$P(\gamma g)=4\frac{\alpha ^2}{\left(1\beta \right)^2}\mathrm{sin}^2\left(\frac{1\beta }{2}\mathrm{\Delta }_\lambda u\right).$$
(58)
When the oscillation length $`\mathrm{}_{\mathrm{osc}}=\frac{2\pi \vartheta }{\mathrm{\Delta }_M}`$ is large with respect to the coherent path distance $`u`$, the weak mixing probability can be further approximated (with $`\mathrm{\Delta }_Mu\vartheta 1`$) by
$$P(\gamma g)\left(\alpha \mathrm{\Delta }_\lambda u\right)^2\left(\mathrm{\Delta }_Mu\right)^2.$$
(59)
* The strong mixing regime in which the mixing is maximal, i.e. when $`\vartheta \pi /4`$, so that the oscillation probability reduces to
$$P(\gamma g)=\mathrm{sin}^2\left(\mathrm{\Delta }_Mu\right)$$
(60)
and the oscillation length to
$$\mathrm{}_{\mathrm{osc}}=\frac{\pi }{\mathrm{\Delta }_M}.$$
(61)
A complete transition between a photon and the light particle is then possible. This can only happen when $`\mathrm{\Delta }_m`$ and $`\mathrm{\Delta }_\lambda `$ have the same sign (see equation (49)). We further note here that the width in $`\beta `$ of the strong mixing region is roughly given by $`\alpha `$ according to equation (49).
### B Mixing in $`D`$ dimensions
#### 1 General result
To compute the oscillation probability in a spacetime with extra–dimensions, we first have to solve (27) which implies the diagonalisation of the matrix (28). We present the explicit and detailed computation of both the eigenvalues and eigenvectors in appendix A. We then compute in appendix B the explicit form of the oscillation probability (see equation (B5))
$$P(\gamma g)=1\left|\underset{i=1}{\overset{N_D}{}}f_{x_i}^2\text{e}^{ix_iu}\right|^2.$$
(62)
Taking into account the fact that $`f_{x_i}^2=1`$, it can be rewritten as
$$P(\gamma g)=2\underset{i,j=1}{\overset{N_D}{}}f_{x_i}^2f_{x_j}^2\mathrm{sin}^2\left[\frac{x_ix_j}{2}u\right],$$
(63)
where the coefficients $`f_{x_j}^2`$ are defined by (see equation (B6))
$$f_{x_j}^2\left[1+\alpha ^2\underset{i=1}{\overset{N_D}{}}\frac{s_i}{\left(y_j\beta _i\right)^2}\right]^1.$$
(64)
The expressions (63) and (64) depend on the eigenvalues $`y_i`$ solutions of the equation (A10). Introducing the notations $`yx/\mathrm{\Delta }_\lambda `$, $`\alpha \mathrm{\Delta }_M/\mathrm{\Delta }_\lambda `$, $`\beta \mathrm{\Delta }_m/\mathrm{\Delta }_\lambda `$ and $`\beta _i\mathrm{\Delta }_m^{(r_i)}/\mathrm{\Delta }_\lambda `$, the eigenvalues equation (A10) can be rewritten as
$$y1=\alpha ^2\underset{i=1}{\overset{N_D}{}}\frac{s_i}{y\beta _i}.$$
(65)
The photon–KK graviton oscillations is then completely described by the set of equations (6365).
Indeed, it is difficult (even impossible if $`n>1`$) to compute analytically the roots of (65). For instance the coefficients $`s_i`$ are not known analytically if $`n>1`$; it is of course possible to compute $`P(\gamma g)`$ numerically, but this is not our purpose here. In the next two sections we derive the oscillation probability in the two cases $`n=1`$ and $`n=2`$ in a range of parameters dictated by the systems where such a mixing may appear. In the next paragraph, we discuss qualitatively the results found there, stressing some new effects due to the presence of a large number of KK states, as well as to the degeneracy of each KK level.
#### 2 Qualitative discussion
We only discuss the cases where $`\alpha `$ is small in comparison to $`\beta `$ as dictated by the physical systems studied in § VIII and § IX. Two different limiting regimes appear, a large radius regime (when $`|\beta |`$ is smaller than unity) and for which there is a significant effect of the extra–dimensions, and a small radius regime (when $`|\beta |`$ is larger than unity) and for which there is only small departure from the usual four dimensional photon mixing.
Let us first discuss the large radius regime where $`|\beta |`$ is smaller than unity. As in four dimensions, according to the respective value of $`\mathrm{\Delta }_\lambda `$ and of the $`\mathrm{\Delta }_m^{(q)}`$’s, two behaviours can appear:
* A strong mixing regime with one given KK state , if there exists a state $`K`$ such that
$$|\mathrm{\Delta }_m^{(K)}\mathrm{\Delta }_\lambda |\sqrt{s_K}\mathrm{\Delta }_M.$$
(66)
This can happen only if $`\beta >0`$, i.e. when plasma effects dominate over the vacuum polarisation. We stress also that since we have assumed along this discussion that $`\alpha `$ is lower than $`\beta `$ there is at most one KK state which can mix strongly with the photon. The total probability will be found to be dominated by a term of the form
$$P(\gamma g)=(1\eta )\mathrm{sin}^2\left[\mathrm{\Delta }_{\mathrm{osc}}^{(K)}u\right]+4\underset{iK}{}\frac{s_i\alpha ^2}{(1\beta _i)^2}\mathrm{sin}^2\left(\frac{1\beta _i}{2}\mathrm{\Delta }_\lambda u\right),$$
(67)
with
$$\mathrm{\Delta }_{\mathrm{osc}}^{(K)}=\sqrt{s_K}\mathrm{\Delta }_M$$
(68)
(see § V and § VI for a detailed derivation). $`\eta `$ is much smaller than unity. As will be shown later, this form accounts for keeping only the dominant part of each $`f_{x_i}^2`$. Depending on the argument of the sines, the probability is either dominated by $`\mathrm{sin}^2\left[\mathrm{\Delta }_{\mathrm{osc}}^{(K)}u\right]`$ or by the correction coming from the modes $`iK`$. This shows a first departure to the four dimensional case due to the degeneracy of the KK level $`K`$; the oscillation length associated with the strong mixing state (labelled by $`K`$) is lowered by a factor $`\sqrt{s_K}`$ which can be very large. Moreover, the width of the region in $`\mathrm{\Delta }_\lambda `$ of strong mixing is larger by a factor $`\sqrt{s_K}`$ than in the usual case \[see below equation (61)\].
An other important difference with the usual four dimensional situation, where the strong mixing regime can only occur when $`\mathrm{\Delta }_\lambda `$ crosses the unique $`\mathrm{\Delta }_m`$ characteristic of the mixing state, we now have more possibilities to be in that regime, where a complete transition between the photon and the graviton is possible. Because of the presence of a KK state $`\mathrm{\Delta }_m^{(q)}`$ in any interval in $`\mathrm{\Delta }_\lambda `$ of typical width $`\mathrm{\Delta }_m`$, only fluctuations of $`\mathrm{\Delta }_\lambda `$ of order $`\mathrm{\Delta }_m`$ can lead to it.
* A weak mixing regime where for all $`q`$, $`|\beta _q1|\alpha `$. The oscillation probability is then given by
$$P(\gamma g)4\underset{i}{}\frac{s_i\alpha ^2}{(1\beta _i)^2}\mathrm{sin}^2\left(\frac{1\beta _i}{2}\mathrm{\Delta }_\lambda u\right).$$
(69)
This contribution is exactly the one that will be intuitively thought of and obtained by summing the individual oscillation probabilities (58) of the photon into each KK state with the mixing angle
$$\mathrm{tan}2\vartheta _q\frac{2\alpha }{1\beta _q}.$$
(70)
There are roughly three contributions to the sum (69) that we estimate as follows.
1. All the states such that $`|\beta _q|1`$ mix with the photon with approximatively the same angle $`\vartheta _q\alpha `$ if we neglect $`\beta _q`$ with respect to unity in (70). The order of magnitude of the probability of oscillations with these states is then
$$P(\gamma g)4𝒩_1\alpha ^2\mathrm{sin}^2\left(\frac{\mathrm{\Delta }_\lambda }{2}u\right).$$
(71)
$`𝒩_1`$ can be estimated by counting the number of modes such that $`\beta _q\beta _{N_1}1`$ with their multiplicity (see appendix C), i.e. $`𝒩_1_{k=1}^{N_1}k^{n1}N_1^n\beta ^{n/2}`$ so that
$$P(\gamma g)\frac{\alpha ^2}{\beta ^{n/2}}\mathrm{sin}^2\left(\frac{\mathrm{\Delta }_\lambda }{2}u\right).$$
(72)
This already shows that the oscillation probability can be greatly enhanced (by a factor $`\beta ^{n/2}`$ with respect to the four dimensional case with the same mixing parameters \[obtained from equation (58) with $`|\beta |1`$\]).
2. All $`\beta _q`$ such that $`|\beta _q|1`$ have a mixing angle roughtly estimated by $`\vartheta _q\alpha /\beta _q`$ and their contribution to the probability is of order
$$P(\gamma g)4\alpha ^2\underset{\beta _q>1}{}\frac{1}{\beta _q^2}\mathrm{sin}^2\left(\frac{\mathrm{\Delta }_\lambda \beta _q}{2}u\right).$$
(73)
This series is difficult to evaluate since the oscillation length is different for each KK state. When $`n3`$ it can be bounded by $`\alpha ^2/\beta ^{n/2}`$ so that this contribution is at most of the same order of magnitude than the previous one.
3. The contribution of the $`\beta _q`$ such that $`\beta _q1`$ which only exists if $`\beta >0`$ is bounded by $`\alpha ^2_{\beta _i1}s_i/(\beta _i1)^2`$. First of all, since $`|\beta _q1|\alpha ^2`$ we are never in a strong mixing regime. Now, we single out $`\beta _K`$, the closest $`\beta _i`$ to unity from which it follows that $`iK`$, $`|\beta _i1|\beta /2`$ and thus the contribution of all the $`\beta _i1`$ for $`iK`$ is bounded by $`(\alpha ^2/\beta ^2)_{\beta _i1,iK}s_i`$. It can be dominated by the contribution of the term $`K`$ given by $`\alpha ^2s_K/|\beta _K1|^2`$ according to the relative value of $`|\beta _K1|`$ in units of $`\beta `$.
In conclusion, the weak mixing case is characterised by an enhancement of the probability by a factor at least $`\beta ^{n/2}`$ due to the fact that the photon couples to a large number of KK states. We further note that when $`\beta <0`$ and $`|\beta |1`$ one can obtain an absolute bound on the oscillation probability of order $`\alpha ^2/\beta ^2`$ for $`n3`$ (when $`n=2`$, this bound is given by $`10𝒬\alpha ^2/\beta ^2`$ \[see appendix (C)\]).
We now turn to the small radius regime where $`|\beta |1`$ and in which the photon mixes preferentially with the zero mode. The probability (63) can be expressed as
$$P(\gamma g)=(1ϵ)P_{4D}(\gamma g)+4\underset{i>1}{}\frac{\alpha ^2}{(1\beta _i)^2}\mathrm{sin}^2\left(\frac{\mathrm{\Delta }_\lambda (1\beta _i)}{2}u\right),$$
(74)
where $`P_{4D}`$ is the oscillation probability for the mixing with the zero mode and is given by (55) and $`ϵ1`$ is the correction of the oscillation probability with this mode coming from the existence of the extra–dimensions. In this case, the lightest massive KK mode is so heavy compared to the photon effective mass that it can barely be excited by the photon. The contribution of the other KK modes can be shown to be bounded by $`𝒪(\alpha ^2\beta ^2)`$. The contribution of the massive KK states is suppressed by a factor $`\beta ^21`$.
Introducing the Compton wavelength, $`\lambda _\gamma `$ say, associated with the effective mass of the photon and defined by
$$\lambda _\gamma |\omega \mathrm{\Delta }_\lambda |^{1/2},$$
(75)
the required condition to be in a large radius regime can be rephrased as $`\lambda _\gamma <R`$, i.e. that the average scale associated with the photon is smaller than the radius of the extra–dimensions. The latter is expected to be of the order of the centimeter for two extra–dimensions.
In the two following sections, we derive these results in details for a five and six dimensional spacetime. Let us stress here than when $`n>3`$ the oscillation probability will strongly depend on the cut–off in energy in which case a more precise knowledge of the whole theory and the exact experimental situation (to take into account decoherence effect) are needed. We emphasize that in the following we have set the cut–off to its maximum value in order to be very general. The computed probability is thus the maximum one and the bounds on the parameters used to derive it are the most stringent.
## V Estimation of the probability in a five dimensional spacetime
In this section, we present the computation of the eigenvalues and of the oscillation probability when $`n=1`$. As it will be seen the computation is easier in this case because the sums in (64-65) can be expressed in terms of circular or hyperbolic functions. Although the case $`n=1`$ is generally regarded as in contradiction with observation (see however e.g. ), this computation will teach us a lot about the mixing with a large number of particles.
In a five dimensional world, the mixing matrix $``$ is explicitely given by
$$_\lambda =\left(\begin{array}{ccccccc}\mathrm{\Delta }_\lambda & \mathrm{\Delta }_M& \mathrm{\Delta }_M& \mathrm{\Delta }_M& \mathrm{\Delta }_M& \mathrm{\Delta }_M& \mathrm{}\\ \mathrm{\Delta }_M& 0& 0& 0& 0& 0& \mathrm{}\\ \mathrm{\Delta }_M& 0& \frac{1}{2R^2\omega }& 0& 0& 0& \mathrm{}\\ \mathrm{\Delta }_M& 0& 0& \frac{1}{2R^2\omega }& 0& 0& \mathrm{}\\ \mathrm{\Delta }_M& 0& 0& 0& \frac{4}{2R^2\omega }& 0& \mathrm{}\\ \mathrm{\Delta }_M& 0& 0& 0& 0& \frac{4}{2R^2\omega }& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\end{array}\right).$$
(76)
This matrix can be compared to the one obtained for neutrino oscillations is spacetime with extra–dimensions (see e.g. ). We see on that example that for $`q0`$ each $`\mathrm{\Delta }_m^{(q)}`$ is twice degenerated so that $`r_1=1`$, $`r_i=2i1`$ and $`s_i=2`$ for $`i>1`$.
The characteristic eigenvalues equation $`\text{det}(xI)=0`$, in that simple case, can be obtained by developing the determinant of order $`2N+2`$ with respect to the last line to get a recursion relation with the determinant of order $`2N`$ and to find the limit of this series. Indeed it leads to the same result that the general equation (A5).
With the notations of the former paragraph, the eigenvalues equation (A5) can now be rewritten after resummation (see 1.217 in) as
$$y1=\frac{\alpha ^2}{\beta }𝒦(y/\beta ),$$
(77)
where the function $`𝒦`$ is defined by
$$𝒦(x)\pi |x|^{1/2}\{\begin{array}{cc}\mathrm{cot}\pi |x|^{1/2}\hfill & (x>0)\hfill \\ \mathrm{coth}\pi |x|^{1/2}\hfill & (x<0).\hfill \end{array}$$
(78)
The oscillation probability is then given by (63) where the coefficients (64) are now reexpressed after resummation as
$$f_{x_j}^2=\left[1+\frac{\alpha ^2}{\beta ^2}(y_j/\beta )\right]^1$$
(79)
where the function $`_{kZ}(xk^2)^2`$ is obtained from (78) as
$$(x)\frac{\pi }{2|x|^{3/2}}\{\begin{array}{cc}\mathrm{cot}\pi \sqrt{x}+\pi \sqrt{x}(1+\mathrm{cot}^2\pi \sqrt{x})\hfill & (x>0)\\ \mathrm{coth}\pi \sqrt{|x|}\pi \sqrt{|x|}(1\mathrm{coth}^2\pi \sqrt{|x|})\hfill & (x<0).\end{array}$$
(80)
We now determine an approximation of the solutions of the eigenvalues equation (77) and then of the oscillation probability in the cases $`\beta <0`$ and $`\beta >0`$ assuming that $`\alpha ^21`$.
### A $`\beta <0`$
#### 1 Eigenvalues
Setting $`\overline{\beta }\beta >0`$, we have $`\beta _j=(j1)^2\overline{\beta }`$ and one can easily sort out that the solutions of equation (77) are such that
$`y_1>1,y_{j+1}]\beta _{j+1},\beta _j[`$ (81)
with
$`y_1`$ solution of $`y1=\pi {\displaystyle \frac{\alpha ^2}{\overline{\beta }}}\sqrt{{\displaystyle \frac{\overline{\beta }}{y}}}\mathrm{coth}\pi \sqrt{{\displaystyle \frac{y}{\overline{\beta }}}},`$ (82)
$`y_{j>1}`$ solutions of $`1y=\pi {\displaystyle \frac{\alpha ^2}{\overline{\beta }}}\sqrt{{\displaystyle \frac{\overline{\beta }}{y}}}\mathrm{cot}\pi \sqrt{{\displaystyle \frac{y}{\overline{\beta }}}}.`$ (83)
In figure 1, we depict the graphical resolution of this equation.
The resolution of equation (77) then splits into the three following cases:
1. $`y_1`$: When $`y_1/\overline{\beta }1`$, (82) reduces to $`y1=\alpha ^2/y`$ so that $`y_11+\alpha ^2`$, and one can then check that the condition $`y_1/\overline{\beta }1`$ is equivalent to $`\overline{\beta }1`$.
When $`y_1/\overline{\beta }1`$, we set $`y_1=1+ϵ`$ with $`ϵ>0`$ and (82) implies that $`0<ϵ=\pi \alpha ^2(1+ϵ)^{1/2}/\overline{\beta }^{1/2}<\pi \alpha ^2/\overline{\beta }^{1/2}`$. Then if $`\alpha ^2/\overline{\beta }^{1/2}1`$ we deduce that $`y_11+\pi \alpha ^2/\overline{\beta }^{1/2}`$ and that the initial condition on $`y_1`$ is equivalent to $`\overline{\beta }1`$.
In conclusion, if $`\alpha ^2/\overline{\beta }^{1/2}1`$,
$$y_1\{\begin{array}{cc}1+\alpha ^2\hfill & (\overline{\beta }1)\\ 1+\pi \frac{\alpha ^2}{\overline{\beta }^{1/2}}\hfill & (\overline{\beta }1).\end{array}$$
(84)
Note that these two solutions can be rewritten under the more compact form
$$y_11+\frac{\alpha ^2}{\beta }𝒦(\beta ^1).$$
(85)
2. $`y_{j+1};j>1`$: Since $`y_{j+1}`$ satisfies
$$(j1)<\sqrt{\frac{y_{j+1}}{\overline{\beta }}}<j$$
we set $`\sqrt{y_{j+1}/\overline{\beta }}(j1)+ϵ_{j+1}`$ with $`0<ϵ_{j+1}<1`$ and equation (83) rewrites as
$`1+\overline{\beta }(j1)^2\left[1+{\displaystyle \frac{ϵ_{j+1}}{j1}}\right]^2={\displaystyle \frac{\pi \alpha ^2}{\overline{\beta }(j1)}}{\displaystyle \frac{\mathrm{cot}\pi ϵ_{j+1}}{1+\frac{ϵ_{j+1}}{j1}}}.`$ (86)
Now, if $`\alpha ^2/\overline{\beta }1`$, the l.h.s. of (86) being larger than unity, it implies that $`\mathrm{cot}\pi ϵ_{j+1}1`$ which thus behaves as $`1/\pi ϵ_{j+1}`$. We can then solve (86) for $`ϵ_{j+1}`$ to get
$$y_j\beta _{j1}\frac{2\alpha ^2}{1\beta _{j1}}.$$
(87)
This expansion is valid whatever the magnitude of $`\overline{\beta }`$ as long as $`\alpha ^2/\overline{\beta }1`$.
3. $`y_2`$: $`y_2`$ is the solution of (83) such that $`0<\sqrt{y_2/\overline{\beta }}<1/2`$. Setting $`zy_2/\overline{\beta }`$, (83) leads to
$`1+\overline{\beta }z=\pi {\displaystyle \frac{\alpha ^2}{\overline{\beta }}}{\displaystyle \frac{\mathrm{cot}\pi \sqrt{z}}{\sqrt{z}}}`$ (88)
with $`0<\sqrt{z}<1/2`$. The l.h.s. of (88) being greater than unity, it implies that, when $`\alpha ^2/\overline{\beta }1`$, $`\mathrm{cot}(\pi \sqrt{z})/\sqrt{z}1`$ and thus behaves as $`1/\pi \sqrt{z}`$. At lowest order (88) then leads to $`z\alpha ^2/\overline{\beta }`$ and then
$$y_2\alpha ^2.$$
(89)
Again, this solution is valid whatever $`\overline{\beta }`$ such that $`\alpha ^2/\overline{\beta }1`$.
4. Summary: When $`\alpha ^21`$, the roots of (77) are well approximated by
$`y_1`$ $``$ $`1+{\displaystyle \frac{\alpha ^2}{\beta }}𝒦(\beta ^1)`$ (90)
$`y_{j>1}`$ $``$ $`\beta _{j1}{\displaystyle \frac{s_{j1}\alpha ^2}{1\beta _{j1}}}`$ (91)
for all $`\beta `$ such that $`\alpha ^2|\beta |`$.
#### 2 Oscillation probability
Assuming that $`\alpha /|\beta |1`$, we can now expand (79) to get the following behaviours of the coefficients $`f_{x_i}^2`$
$`f_{x_1}^2`$ $``$ $`1{\displaystyle \frac{\alpha ^2}{\beta ^2}}(\beta ^1),`$ (92)
$`f_{x_{j>1}}^2`$ $``$ $`{\displaystyle \frac{s_{j1}}{(1\beta _{j1})^2}}\alpha ^2.`$ (93)
One can easily check that, at this order, $`f_{x_i}^2=1`$. Using the form (63) of the oscillation probability we deduce that
$$P(\gamma g)4\alpha ^2\underset{j1}{}\frac{s_j}{(1\beta _j)^2}\mathrm{sin}^2\left(\frac{1\beta _j}{2}\mathrm{\Delta }_\lambda u\right),$$
(95)
as announced in (69). It can be checked that the dominant contribution to the probability comes from the terms $`f_{x_1}^2f_{x_j}^2`$ in (63).
It is worth noting that the cut–off of the theory does not enter the result, due to the fact that in this special case all the sums are converging. Since $`\beta _j0`$ and $`\alpha 1`$, we are always in the weak mixing limit and the oscillation probability is well approximated by the sum of all the individual oscillation probabilities. The individual oscillation lengths are given by
$$\mathrm{}_{\mathrm{osc}}^{(j)}=\frac{2\pi }{\mathrm{\Delta }_\lambda }\frac{1}{1\beta _j}<\frac{2\pi }{\mathrm{\Delta }_\lambda }=\mathrm{}_{\mathrm{osc}}^{(1)}.$$
### B $`\beta >0`$
#### 1 Eigenvalues
Since $`\beta _j=(j1)^2\beta `$, one can easily show that the roots of (77) are such that
$$y_1<0,\beta _i<y_{j+1}<\beta _{j+1}$$
(96)
with
$`y_1`$ solution of $`1y=\pi {\displaystyle \frac{\alpha ^2}{\beta }}\sqrt{{\displaystyle \frac{\beta }{y}}}\mathrm{coth}\pi \sqrt{{\displaystyle \frac{y}{\beta }}},`$ (97)
$`y_{j>1}`$ solutions of $`y1=\pi {\displaystyle \frac{\alpha ^2}{\beta }}\sqrt{{\displaystyle \frac{\beta }{y}}}\mathrm{cot}\pi \sqrt{{\displaystyle \frac{y}{\beta }}}.`$ (98)
In figure 2 we depict the graphic resolution of this equation. We introduce $`K`$ the index of the closest $`\beta _i`$ to unity. Contrary to the previous case, the discussion has to be split in four steps:
1. $`y_1`$: When $`y_1/\beta 1`$, (97) implies that $`1y_1\alpha ^2/y_1`$ so that $`y_1\alpha ^2`$ and the initial condition reduces to $`\alpha ^2/\beta 1`$. Thus when $`\alpha ^2/\beta 1`$, whatever the magnitude of $`\beta `$,
$$y_1\alpha ^2.$$
(99)
2. $`y_j,1<j<K`$: We set $`\sqrt{y_j/\beta }=(j1)ϵ_j`$ with $`1>ϵ_j>0`$, so that (98) rewrites as
$$1y_j=\frac{\pi \alpha ^2}{\beta }\frac{\mathrm{cot}\pi ϵ_j}{(j1)ϵ_j}.$$
(100)
Since the l.h.s. of (100) is positive we deduce that $`ϵ_j<1/2`$. Now, taking into account the fact that $`1](\beta _K+\beta _{K1})/2,(\beta _K+\beta _{K+1})/2[`$, we deduce that $`1y_j>1y_{K1}>(\beta _K\beta _{K1})/2=\beta (K3/2)/2`$ and thus (since $`K2`$) that $`1y_j\beta /4`$. If $`\alpha ^2/\beta ^21`$, then (100) implies that $`\mathrm{cot}\pi ϵ_j1/\pi ϵ_j1`$ from which we deduce $`ϵ_j`$ and then
$$y_{1<j<K}\beta (j1)^2\frac{2\alpha ^2}{1\beta (j1)^2}.$$
(101)
3. $`y_j,j>K+1`$: The argument follows the same lines as the previous one but we now set $`\sqrt{y_j/\beta }=(j2)+ϵ_j`$ with $`1>ϵ_j>0`$. We can now deduce from $`y_j1y_{K+2}1`$ that $`y_j1\beta /2`$ and then that if $`\alpha ^2/\beta ^21`$,
$$y_{j>K+1}\beta (j2)^2+\frac{2\alpha ^2}{\beta (j2)^21}.$$
(102)
4. $`y_K,y_{K+1}`$: If $`K>1`$, we set $`y_K=\beta _K(1+ϵ)`$ and $`y_{K+1}=\beta _K(1+ϵ^{})`$, we can use the property of equation (98) to deduce, as before, that
$$0<ϵ<\frac{K3/4}{(K1)^2}\mathrm{and}0<ϵ^{}<\frac{K5/4}{(K1)^2}$$
(103)
and then if $`\beta `$ is small, we can conclude that $`ϵ`$ and $`ϵ^{}`$ are small compared to unity. Setting $`\delta \beta _K1`$, $`ϵ`$ and $`ϵ^{}`$ are solution of (98) which reduces to
$$\delta +ϵ\frac{2\alpha ^2}{ϵ}2ϵ\delta \pm \sqrt{\delta ^2+8\alpha ^2}.$$
(104)
Now, if $`\delta 2\sqrt{2}\alpha `$, we deduce that
$$y_K,y_{K+1}\{\beta _K+\frac{2\alpha ^2}{\beta _K1},1\frac{2\alpha ^2}{\beta _K1}\},$$
(105)
depending on the sign of $`\delta `$ and with the constraint $`y_K<y_{K+1}`$. On the other hand, if $`\delta 2\sqrt{2}\alpha `$,
$$y_K1+\alpha \sqrt{2},y_{K+1}1\alpha \sqrt{2}.$$
(106)
Note that if $`K=1`$, then $`\beta 2`$ and the above discussion is still valid, but we just have the two classes of solutions $`y_1`$ and $`y_{j>1}`$.
#### 2 Oscillation probability
Assuming that $`\alpha /\beta 1`$, we can expand (79) to get the following forms for the coefficients $`f_{x_i}^2`$. Assuming that $`\beta _K1<0`$,
$`f_{x_{1j<K}}^2`$ $``$ $`{\displaystyle \frac{s_j}{(1\beta _j)^2}}\alpha ^2`$ (107)
$`f_{x_K}^2`$ $``$ $`\{\begin{array}{cc}\frac{1}{2}\hfill & (|\beta _K1|2\sqrt{2}\alpha )\\ \frac{s_K}{(1\beta _K)^2}\alpha ^2\hfill & (|\beta _K1|2\sqrt{2}\alpha )\end{array}`$ (110)
$`f_{x_{K+1}}^2`$ $``$ $`\{\begin{array}{cc}\frac{1}{2}\hfill & (|\beta _K1|2\sqrt{2}\alpha )\\ 1\frac{\alpha ^2}{\beta ^2}(\beta ^1)\hfill & (|\beta _K1|2\sqrt{2}\alpha )\end{array}`$ (113)
$`f_{x_{j>K+1}}^2`$ $``$ $`{\displaystyle \frac{s_{j1}}{(1\beta _{j1})^2}}\alpha ^2.`$ (114)
When $`\beta _K1>0`$, $`f_{x_{K1}}^2`$ is then given by (113) and $`f_{x_{K+1}}^2`$ is of the form (114).
It can be checked that, at lowest order $`f_{x_i}^2=1`$. One can further check that when $`|\beta _K1|2\sqrt{2}\alpha `$ the first correction to (110113) is $`\alpha ^2(\beta ^1)/2\beta ^2`$ and that $`f_{x_i}^2=1`$ is also satisfied (which indeed has to be order by order). Now, the oscillation probability (63) reduces to
$$P(\gamma g)\{\begin{array}{cc}(1ϵ)\mathrm{sin}^2\left[\mathrm{\Delta }_Ms_K^{1/2}u\right]+4\alpha ^2_{iK}\frac{s_i}{(1\beta _i)^2}\mathrm{sin}^2\left[\frac{1\beta _i}{2}\mathrm{\Delta }_\lambda u\right]\hfill & (|\beta _K1|2\sqrt{2}\alpha )\\ 4\alpha ^2_{i1}\frac{s_i}{(1\beta _i)^2}\mathrm{sin}^2\left[\frac{1\beta _i}{2}\mathrm{\Delta }_\lambda u\right]\hfill & (|\beta _K1|2\sqrt{2}\alpha ),\end{array}$$
(115)
with $`\eta \alpha ^2(\beta ^1)/\beta ^2`$ as announced in (69). Thus, when $`|\beta _K1|\alpha `$, we are in a weak mixing regime and the probability is obtained by summing over all the individual probabilities. Otherwise, we are in a regime of strong mixing with the state $`K`$, and the oscillation length with this state is now given by
$$\mathrm{}_{\mathrm{osc}}=\frac{\pi }{\mathrm{\Delta }_Ms_K^{1/2}}.$$
### C Summary and discussion
In the limit where $`\alpha ^21`$, we have estimated the oscillation probability (63) to be
* $`\beta <0`$:
$$P(\gamma g)4\alpha ^2\underset{i1}{}\frac{s_i}{(1\beta _i)^2}\mathrm{sin}^2\left[\frac{1\beta _i}{2}\mathrm{\Delta }_\lambda u\right],$$
* $`\beta >0`$:
$$P(\gamma g)\{\begin{array}{cc}I\mathrm{sin}^2\left[\mathrm{\Delta }_Ms_K^{1/2}u\right]+4\alpha ^2_{iK}\frac{s_i}{(1\beta _i)^2}\mathrm{sin}^2\left[\frac{1\beta _i}{2}\mathrm{\Delta }_\lambda u\right]\hfill & (|\beta _K1|2\sqrt{2}\alpha )\\ 4\alpha ^2_{i1}\frac{s_i}{(1\beta _i)^2}\mathrm{sin}^2\left[\frac{1\beta _i}{2}\mathrm{\Delta }_\lambda u\right]\hfill & (|\beta _K1|2\sqrt{2}\alpha ^2),\end{array}$$
with $`I1\alpha ^2(\beta ^1)/\beta ^2`$, as announced in (67) and (69). When $`|\beta _j1|2\sqrt{2}\alpha ^2`$ for all $`j`$, it is interesting to study the limit where $`u\mathrm{}_{\mathrm{osc}}`$ in which case we can assume that the sines can be replaced by their average value to get the two limiting behaviours:
$`P(\gamma g)`$ $``$ $`2\alpha ^2\left(1+{\displaystyle \frac{\pi ^4}{45\beta ^2}}\right)(|\beta |1),`$ (116)
$`P(\gamma g)`$ $``$ $`\pi {\displaystyle \frac{\alpha ^2}{\sqrt{|\beta |}}}(|\beta |1).`$ (117)
As explained in § IV, we respectively see on (116-117) the small and large radius regimes where the extra–dimensions either have no effect (116) or enhance (117) the probability. We also find the regime of strong mixing with the $`K^{\mathrm{th}}`$ KK graviton and the effect on the oscillation length as discussed below equation (68).
## VI Estimation of the probability in a six dimensional spacetime
Let us now turn to the physically more interesting case of $`n=2`$ extra–dimensions. Now, equation (65) cannot be solved exactly in general, but its solutions can be well approximated when the coupling between the photon and the graviton (this coupling is measured by $`\mathrm{\Delta }_M`$ in the matrix (28)) is small enough compared to the typical mass parameters of the mixing particles \[i.e. the diagonal terms in (28)\]. We follow the same lines as in the previous section.
### A $`1\beta >0`$
#### 1 Eigenvalues
In that case, the plasma effects dominate over the vacuum polarisation in $`\mathrm{\Delta }_\lambda `$; all the $`\beta _i`$ are positive and one can easily show that the roots of (65) are such that
$`y_1<0,y_i]\beta _{i1},\beta _i[,y_{N_D+1}>\beta _{N_D}.`$ (118)
We introduce $`k(i)`$ the index such that $`\beta _{k(i)}`$ is the closest $`\beta _i`$ to $`y_i`$. From (118) one has that $`k(i)\{i1,i\}`$. The eigenvalue equation (65) for the root $`y_i`$ can be rewritten as
$$y_i1=\alpha ^2\frac{s_k}{y_i\beta _k}+_k(y_i),$$
(119)
where $`_k`$ is defined as
$$_k(y)\alpha ^2\underset{jk}{}\frac{s_j}{y\beta _j}.$$
(120)
$``$ with no subscript denotes the function defined by the sum (120) taken over all indices $`j`$ from $`j=1`$ to $`j=N_D`$. To finish, we introduce the index $`K`$ such that $`\beta _K`$ is the closest $`\beta _i`$ to unity and then,
$$iK|\beta _i1|\frac{\beta }{2}.$$
(121)
As in the five dimensional case, the determination of the roots $`y_i`$ has to be split in the three following cases:
* $`iK1`$: We first show that $`k(i)=i`$. For that purpose, we consider the function $`(y)`$ defined by
$$(y)1y+(y).$$
(122)
This function is strictly decreasing on $`]\beta _{i1},\beta _i[`$ ($``$ vanishes only once in this interval in $`y=y_i`$). Showing that $`\left(\frac{\beta _i+\beta _{i1}}{2}\right)0`$, is then enough to prove that $`k(i)=i`$. Since one has
$$1\frac{\beta _i+\beta _{i1}}{2}\frac{\beta }{2}\text{ and }k\left|\frac{\beta _i+\beta _{i1}}{2}\beta _k\right|\frac{\beta }{2},$$
(123)
using (C19), one obtains
$$\left|\left(\frac{\beta _i+\beta _{i1}}{2}\right)\right|𝒬\frac{\alpha ^2}{\beta ^2}\text{sup}(𝒬^{},\sqrt{\beta y})𝒬𝒬^{}\frac{\alpha ^2}{\beta ^2}\text{for }y1\text{ and }\beta 1.$$
(124)
(the constants $`𝒬`$ and $`𝒬^{}`$ are defined in equation (C19) of appendix C). Comparing (123) and (124) one sees that for $`\alpha `$ smaller enough than $`\beta `$ (namely $`2𝒬𝒬^{}\frac{\alpha ^2}{\beta ^3}<1)`$, one has $`\left(\frac{\beta _i+\beta _{i1}}{2}\right)>0`$ and then that $`k(i)=i`$.
We will assume in the following the slightly stronger constraint
$$10𝒬𝒬^{}\frac{\alpha ^2}{\beta ^3}<1.$$
(125)
Now, we set $`y_i=\beta _iϵ_i`$, with $`ϵ_i>0`$. Equation (119) can be rewritten as an equation for $`ϵ_i`$, with $`_i_i(y_i)`$,
$$\frac{ϵ_i^2}{\beta _i1_i}ϵ_i\frac{\alpha ^2s_i}{\beta _i1_i}=0,$$
(126)
the positive solution of which is given to leading order by
$$ϵ_i\frac{\alpha ^2s_i}{1\beta _i},$$
(127)
when $`\alpha `$ and $`\beta `$ verify the constraint (125) to establish this we have used (C19) and (C21) and showed that (125) leads to $`|\beta _i1|_i`$ and $`|\beta _i1|2\alpha \sqrt{s_i}`$ which in turn leads to the expression (127). In the rest of this section $`ab`$ means that $`a>10b`$ which we assumed to be enough to neglect b with respect a.. The eigenvalues are then given at leading order by
$$y_i\beta _i\alpha ^2\frac{s_i}{1\beta _i}.$$
(128)
* $`iK+2`$: Using a similar line of reasoning as in the previous case, one can show that $`k(i)=i1`$ and then that $`y_i`$ is given at dominant order by (for $`\alpha `$ and $`\beta `$ verifying (125))
$$y_i\beta _{i1}+\alpha ^2\frac{s_{i1}}{\beta _{i1}1}.$$
(129)
* $`i=K,K+1`$: We first estimate the root $`y_K`$. We assume that $`1[\beta _K,\beta _{K+1}[`$ (similar conclusions can be obtained when $`1[\beta _{K1},\beta _K[`$), we have
$$1\frac{\beta _K+\beta _{K1}}{2}>\frac{\beta }{2}.$$
(130)
As in the previous case, this is enough to show that $`y_K[\beta _{K1}+\frac{\beta _K\beta _{K1}}{2},\beta _K[`$ <sup>§</sup><sup>§</sup>§here we use (125) again.. We set $`y_K=\beta _Kϵ_K`$ and $`_K_K(y_K)`$ with $`ϵ_K>0`$ solution of
$$\frac{ϵ_K^2}{\beta _K1_K}ϵ_K\frac{\alpha ^2s_K}{\beta _K1_K}=0,$$
(131)
the positive root of which is
$$ϵ_K=\frac{1\beta _K+_K}{2}\left(1+\sqrt{1+\frac{4\alpha ^2s_K}{(\beta _K1_K)^2}}\right).$$
(132)
For $`\alpha `$ smaller enough than $`\beta `$ One has here to impose a slightly stronger condition than the previous one: namely $`4𝒬𝒬^{}\frac{\alpha }{\beta ^2}<1`$ in order to be able to neglect $`_K`$ with respect to $`2\alpha \sqrt{s_K}`$. This also insures that the expressions (133) are valid. , one can consider the two limiting regimes (we will not consider here the intermediate case, in order to simplify the discussion)
$$ϵ_K\{\begin{array}{ccc}\alpha \sqrt{s_K}\hfill & \text{if}& |\beta _K1|2\alpha \sqrt{s_K}\hfill \\ \frac{\alpha ^2s_K}{1\beta _K}\hfill & \text{if}& |\beta _K1|2\alpha \sqrt{s_K}.\hfill \end{array}$$
(133)
Let us now turn to the evaluation of the root $`y_{K+1}`$. The discussion mimics the previous one. Assuming that 1 is not too close to $`\frac{\beta _K+\beta _{K+1}}{2}`$ namely $`\left|1\frac{\beta _K+\beta _{K+1}}{2}\right|>\frac{\beta }{10}`$. When this is not the case, one obtains the same results as in (134) when $`|\beta _K1|2\alpha \sqrt{s_K}`$., one can show that $`y_{K+1}]\beta _K,\beta _K+\frac{\beta _{K+1}\beta _K}{2}[`$ under the condition (125). Then we write $`y_{K+1}=\beta _Kϵ_{K+1}`$, with $`ϵ_{K+1}<0`$. $`ϵ_{K+1}`$ is solution of equation (131) with $`_K_K(y_{K+1})`$. One considers (under the condition on $`\alpha `$ and $`\beta `$ of footnote $``$VI A 1) the two limiting regimes
$$y_{K+1}\{\begin{array}{ccc}\beta _K+\alpha \sqrt{s_K}\hfill & \text{for}& |\beta _K1|2\alpha \sqrt{s_K},\hfill \\ 1+(1)\hfill & \text{for}& |\beta _K1|2\alpha \sqrt{s_K}.\hfill \end{array}$$
(134)
#### 2 Oscillation probability
We now need to expand the coefficients (64) in order to estimate the probability (63).
* For $`iK1`$, we have from equation (64) and (128)
$$f_{x_i}^2=\left[1+\frac{(1\beta _i)^2}{\alpha ^2s_i}+𝒢_i\right]^1$$
(135)
with
$$𝒢_i𝒢_i(y_i)\alpha ^2\underset{k=1,ki}{\overset{N_D}{}}\frac{s_k}{\left(y_i\beta _k\right)^2}.$$
(136)
Using (C20), one can then show that under the condition (125)
$$\frac{(1\beta _i)^2}{\alpha ^2s_i}\mathrm{max}(1,𝒢_i)$$
(137)
so that, at dominant order,
$$f_{x_i}^2\alpha ^2\frac{s_i}{(1\beta _i)^2}.$$
(138)
* For $`iK+2`$, we find in a similar way (and under the same condition)
$$f_{x_i}^2\alpha ^2\frac{s_{i1}}{(1\beta _{i1})^2}.$$
(139)
* For $`i=K,K+1`$, we distinguish the two regimes $`|\beta _K1|2\alpha \sqrt{s_K}`$ and $`|\beta _K1|2\alpha \sqrt{s_K}`$.
Assuming that $`|\beta _K1|2\alpha \sqrt{s_K}`$, we find <sup>\**</sup><sup>\**</sup>\**Using (C20), we show that this expansion is valid under the condition of the footnote ($``$VI A 1). the dominant contribution to $`f_{x_K}`$ and $`f_{x_{K+1}}`$ to be
$$f_{x_K}^2\alpha ^2\frac{s_K}{(1\beta _K)^2}\text{and}f_{x_{K+1}}^21.$$
(140)
When $`|\beta _K1|2\alpha \sqrt{s_K}`$, the dominant contribution to $`f_{x_K}^2`$ and $`f_{x_{K+1}}^2`$ are
$$f_{x_K}^2\frac{1}{2}\text{and}f_{x_{K+1}}^2\frac{1}{2}.$$
(141)
Now, inserting these results in (63), we find that the oscillation probability is given to dominant order, for $`|\beta _K1|2\alpha \sqrt{s_K}`$, by
$$P(\gamma g)4\underset{iK+1}{}f_{x_i}^2\mathrm{sin}^2\left[\frac{x_ix_{K+1}}{2}u\right]4\underset{i=1}{\overset{i=N_D}{}}\alpha ^2\frac{s_i}{(1\beta _i)^2}\mathrm{sin}^2\left[\frac{1\beta _i}{2}\mathrm{\Delta }_\lambda u\right]$$
(142)
This expression, as announced, is analogous to (69) and corresponds to the case when no KK state mixes strongly with the photon.
Now, when $`|\beta _K1|2\alpha \sqrt{s_K}`$, one has
$`P(\gamma g)`$ $``$ $`(1\eta )\mathrm{sin}^2\left[{\displaystyle \frac{x_Kx_{K+1}}{2}}u\right]+2{\displaystyle \underset{P=K,K+1}{}}{\displaystyle \underset{iP}{}}f_{x_i}^2\mathrm{sin}^2\left[{\displaystyle \frac{x_ix_P}{2}}u\right]`$ (143)
which can be rewritten as
$$P(\gamma g)(1\eta )\mathrm{sin}^2[\mathrm{\Delta }_M\sqrt{s_K}u]+4\underset{iK}{}\alpha ^2\frac{s_i}{(1\beta _i)^2}\mathrm{sin}^2\left[\frac{(1\beta _i)}{2}\mathrm{\Delta }_\lambda u\right]$$
(144)
with $`\eta _{iK}\alpha ^2s_i/(1\beta _i)^21`$. As in the five dimensional case it was obtained by imposing that $`f_{x_i}^2=1`$ order by order. This corresponds to the case where one KK state mixes strongly with the photon and again the oscillation probability is found to be equivalent to (67).
### B $`1\beta <0`$
We now consider the case where $`\beta <0`$ (i.e. when the vacuum contribution dominates over the plasma effects in $`\mathrm{\Delta }_\lambda `$). It is easy to see graphically (see figure 1) that the $`N_D+1`$ solutions of the eigenvalue equation (65) are such that
$$y_1>1,y_{i+1}]\beta _{i+1},\beta _i[,y_{N_D+1}<\beta _{N_D}.$$
(145)
#### 1 Eigenvalues
We do not detail the computation of the eigenvalues since it is similar to the former case. It is even simpler since now we do not have to single out the mode $`K`$ (look for instance to the five dimensional case § V A).
We have only to assume the less stringent constraint than (125), $`10\frac{\alpha ^2}{\beta ^2}𝒬𝒬^{}<1`$, in order for the following expansions to be valid. Under this condition one can show that
$`y_1`$ $``$ $`1+\alpha ^2{\displaystyle \underset{i}{}}{\displaystyle \frac{s_i}{1\beta _i}},`$ (146)
$`y_{j>1}`$ $``$ $`\beta _{j1}{\displaystyle \frac{s_{j1}}{1\beta _{j1}}}\alpha ^2.`$ (147)
#### 2 Oscillation Probability
The coefficients (64) are then given at leading order by
$`f_{x_1}^2`$ $``$ $`1,`$ (148)
$`f_{x_{j>1}}^2`$ $``$ $`\alpha ^2{\displaystyle \frac{s_{j1}}{(1\beta _{j1})^2}}.`$ (149)
The oscillation probability (63) can be expanded as
$$P(\gamma g)4f_{x_1}^2\underset{j>1}{}f_{x_j}^2\mathrm{sin}^2\left[\frac{x_1x_j}{2}u\right]$$
(150)
where we have neglected higher order terms. This can be rewritten as
$$P(\gamma g)4\alpha ^2\underset{j}{}\frac{s_j}{(1\beta _j)^2}\mathrm{sin}^2\left[\frac{1\beta _j}{2}\mathrm{\Delta }_\lambda u\right]$$
(151)
which is again analogous to (69).
### C $`|\beta |>1`$
We introduce the new parameters $`\overline{\alpha }\mathrm{\Delta }_M/\mathrm{\Delta }_m<0`$ and $`\overline{\beta }_i\mathrm{\Delta }_m^{(r_i)}/\mathrm{\Delta }_m=\stackrel{}{p}_{(r_i)}^2`$, $`\overline{\gamma }\mathrm{\Delta }_\lambda /\mathrm{\Delta }_m`$ and $`zx/\mathrm{\Delta }_m`$. The eigenvalue equation (A10) can be rewritten as
$$z\overline{\gamma }=\overline{\alpha }^2\underset{i=1}{\overset{N_D}{}}\frac{s_i}{(z\overline{\beta }_i)}=\frac{\overline{\alpha }^2}{z}+\overline{\alpha }^2\underset{i=2}{\overset{N_D}{}}\frac{s_i}{(z\overline{\beta }_i)}.$$
(152)
One sees easily graphically that this equation admits one negative root $`z_1`$ and that the other $`z_i`$ ($`i3`$) verify $`z_i]\overline{\beta }_{i1},\overline{\beta }_i[`$. We discuss here only the case where $`\overline{\gamma }`$ is closer to $`0`$ than to $`1`$ (we assume that $`\overline{\gamma }<1/4)`$. The discussion is very similar to the previous cases. Under the condition $`10\overline{\alpha }^2𝒬𝒬^{}<1`$, one finds that the roots $`z_i`$ with $`i2`$ are given by
$$z_i\overline{\beta }_{i1}+\frac{\overline{\alpha }^2s_{i1}}{\overline{\beta }_{i1}\overline{\gamma }}\text{so that}f_{x_i}^2\frac{\overline{\alpha }^2s_{i1}}{(\overline{\beta }_{i1}\overline{\gamma })^2}.$$
(153)
Under the condition $`4\overline{\alpha }𝒬𝒬^{}<1`$, the two roots $`z_1`$ and $`z_2`$ are given, for $`\overline{\gamma }>0`$ by<sup>††</sup><sup>††</sup>††For $`\overline{\gamma }<0`$ the results are similar; one has only to exchange the expressions of $`z_1`$ and $`z_2`$ in the case $`|\overline{\gamma }|2|\overline{\alpha }|`$. by
$$\{\begin{array}{c}z_1\overline{\alpha }\hfill \\ z_2\overline{\alpha }\hfill \end{array}\text{if}\overline{\gamma }2|\overline{\alpha }|\text{and by}\{\begin{array}{c}z_1\frac{\overline{\alpha }^2}{\overline{\gamma }}\hfill \\ z_2\overline{\gamma }+(\overline{\gamma })\hfill \end{array}\text{if}\overline{\gamma }2|\overline{\alpha }|$$
(154)
from which we deduce that the coefficients are given either by
$$\{\begin{array}{c}f_{x_1}^2\frac{1}{2}\hfill \\ f_{x_2}^2\frac{1}{2}\hfill \end{array}\text{if}\overline{\gamma }2|\overline{\alpha }|\text{or by}\{\begin{array}{c}f_{x_1}^21\hfill \\ f_{x_2}^2\frac{\overline{\alpha }^2}{\overline{\gamma }^2}\hfill \end{array}\text{if}\overline{\gamma }2|\overline{\alpha }|.$$
(155)
When $`\overline{\gamma }2|\overline{\alpha }|`$, the oscillation probability is given by
$$P(\gamma g)(1\stackrel{~}{\eta })\mathrm{sin}^2(\mathrm{\Delta }_Mu)+4\underset{i2}{}\frac{\alpha ^2s_i}{\beta _i^2}\mathrm{sin}^2\left[\frac{\beta _i}{2}\mathrm{\Delta }_\lambda u\right]$$
(156)
and when $`\overline{\gamma }2|\overline{\alpha }|`$ by
$$P(\gamma g)4\underset{i1}{}\frac{\alpha ^2s_i}{(1\beta _i)^2}\mathrm{sin}^2\left[\frac{1\beta _i}{2}\mathrm{\Delta }_\lambda u\right].$$
(157)
The small coefficient $`\stackrel{~}{\eta }_{i2}s_i\alpha ^2/\beta _i^2`$ is obtained as in (143).
### D Summary
In the limit where $`\alpha ^2<1`$ we have estimated the oscillation probability (63) for a six dimensional spacetime to be
* $`1\beta >0`$: For $`|\beta _K1|2\alpha \sqrt{s_K}`$,
$$P(\gamma g)4\underset{i=1}{\overset{i=N_D}{}}\alpha ^2\frac{s_i}{(1\beta _i)^2}\mathrm{sin}^2\left[\frac{1\beta _i}{2}\mathrm{\Delta }_\lambda u\right]$$
(158)
and for $`|\beta _K1|2\alpha \sqrt{s_K}`$
$$P(\gamma g)(1\eta )\mathrm{sin}^2[\mathrm{\Delta }_M\sqrt{s_K}u]+4\underset{iK}{}\alpha ^2\frac{s_i}{(1\beta _i)^2}\mathrm{sin}^2\left[\frac{(1\beta _i}{2}\mathrm{\Delta }_\lambda u\right],$$
(159)
these results being valid as long as $`4𝒬𝒬^{}\frac{\alpha }{\beta ^2}<1`$ and the coefficient $`\eta `$ being defined in equation (144).
* $`1\beta <0`$:
$$P(\gamma g)4\alpha ^2\underset{j}{}\frac{s_j}{(1\beta _j)^2}\mathrm{sin}^2\left[\frac{1\beta _j}{2}\mathrm{\Delta }_\lambda u\right],$$
(160)
valid if $`10\frac{\alpha ^2}{\beta ^2}𝒬𝒬^{}<1`$.
* $`|\beta |>1`$:
$`P(\gamma g)`$ $``$ $`(1\stackrel{~}{\eta })\mathrm{sin}^2(\mathrm{\Delta }_Mu)+4{\displaystyle \underset{i2}{}}{\displaystyle \frac{\alpha ^2s_i}{\beta _i^2}}\mathrm{sin}^2\left[{\displaystyle \frac{\beta _i}{2}}\mathrm{\Delta }_\lambda u\right]`$ (161)
$``$ $`4{\displaystyle \underset{i1}{}}{\displaystyle \frac{\alpha ^2s_i}{(1\beta _i)^2}}\mathrm{sin}^2\left[{\displaystyle \frac{1\beta _i}{2}}\mathrm{\Delta }_\lambda u\right]`$ (162)
respectively for $`\overline{\gamma }2|\overline{\alpha }|`$ and for $`\overline{\gamma }2|\overline{\alpha }|`$, the result being valid if $`10\overline{\alpha }^2𝒬𝒬^{}<1`$ and the small coefficient $`\stackrel{~}{\eta }`$ is defined in (156).
## VII Mixing in an inhomogeneous field
In all the previous sections, we have assumed that the magnetic field was homogeneous. This is however a very crude approximation for most of the realistic physical systems. In this section we first extand our analysis to inhomogeneous magnetic fields and give some implications of the inhomogeneity of the external field.
### A Computation of the oscillation probability
Following , we rewrite the equation of evolution (27) as a Schrödinger equation
$$i_u\stackrel{}{𝒱}=\left(_0+_1\right)\stackrel{}{𝒱},$$
(163)
where we set $`\stackrel{}{𝒱}(A,G^{(0)},\mathrm{},G^{(N)})`$. The two matrices $`_0`$ and $`_1`$ are respectively defined by
$$_0(u)\omega +\left(\begin{array}{ccccc}\mathrm{\Delta }_\lambda & & & & \\ & 0& & & \\ & & \mathrm{\Delta }_m^{(1)}& & \\ & & & \mathrm{}& \\ & & & & \mathrm{\Delta }_m^{(N)}\end{array}\right)$$
(164)
and
$$_1(u)\left(\begin{array}{cccc}0& \mathrm{\Delta }_M& \mathrm{}& \mathrm{\Delta }_M\\ \mathrm{\Delta }_M& & & \\ \mathrm{}& & & \\ \mathrm{\Delta }_M& & & \end{array}\right).$$
(165)
We assume that $`_1`$ is a perturbation compared to $`_0`$. This approximation is equivalent to saying that $`\mathrm{\Delta }_M/\mathrm{\Delta }_\lambda `$ and $`\mathrm{\Delta }_M/\mathrm{\Delta }_m`$ are small compared to unity, i.e. that $`\alpha 1`$ and $`\alpha /\beta 1`$. When $`H_0`$ is inhomogeneous only $`\mathrm{\Delta }_M`$ and $`\mathrm{\Delta }_\lambda `$ depend on $`u`$ while all the $`\mathrm{\Delta }_m^{(q)}`$ are constant.
We first solve (163) at zeroth order, i.e. by neglecting $`_1`$ with respect to $`_0`$, as
$$\stackrel{}{𝒱}^{(0)}(u)=U(u)\stackrel{}{𝒱}^{(0)}(0),$$
(166)
where the evolution operator $`U`$ is defined by
$$U(u)\mathrm{exp}i_0^u_0(u^{})𝑑u^{}.$$
(167)
Note that at this order there is no mixing effect since $`_0`$ is diagonal.
The general solution of (163) is obtained by shifting to the “interaction representation” where $`\stackrel{}{𝒱}_{\mathrm{int}}U^{}\stackrel{}{𝒱}`$ so that (163) can be rewritten as
$$i_u\stackrel{}{𝒱}_{\mathrm{int}}=_{\mathrm{int}}\stackrel{}{𝒱}_{\mathrm{int}}$$
(168)
with $`_{\mathrm{int}}U^{}_1U`$. This equation can be solved iteratively by setting $`\stackrel{}{𝒱}_{\mathrm{int}}=\stackrel{}{𝒱}_{\mathrm{int}}^{(k)}`$ with
$$\stackrel{}{𝒱}_{\mathrm{int}}^{(n+1)}=i_0^u𝑑u^{}_{\mathrm{int}}(u^{})\stackrel{}{𝒱}_{\mathrm{int}}^{(n)}(u^{}),$$
(169)
with the initial condition $`\stackrel{}{𝒱}_{\mathrm{int}}^{(0)}\stackrel{}{𝒱}(0)`$.
Introducing the basis $`\{\stackrel{}{𝒜},\stackrel{}{𝒢}_q\}`$ with $`\stackrel{}{𝒜}(1,0,\mathrm{},0)`$ and $`\stackrel{}{𝒢}_q`$ being the state of the $`q^{\mathrm{th}}`$ graviton, $`𝒢_q^i=\delta _q^i`$ for $`i\{1,\mathrm{},N+2\}`$ and starting with an initial state describing a pure photon, i.e. $`\stackrel{}{𝒱}(0)=A(0)\stackrel{}{𝒜}`$, we obtain
$$\stackrel{}{𝒱}_{\mathrm{int}}^{(1)}(u)=i_0^u𝑑z\mathrm{\Delta }_M(z)\underset{q}{}\text{e}^{i_0^z(\mathrm{\Delta }_m^{(q)}\mathrm{\Delta }_\lambda (y))𝑑y}A(0)\stackrel{}{𝒢}_q,$$
(170)
where we have used that $`_0\stackrel{}{𝒜}=\mathrm{\Delta }_\lambda \stackrel{}{𝒜}`$, $`_1\stackrel{}{𝒜}=\mathrm{\Delta }_M_q\stackrel{}{𝒢}_q`$ and $`_1\stackrel{}{𝒢}_q=\mathrm{\Delta }_m^{(q)}\stackrel{}{𝒢}_q`$. If we restrict to the first iteration, the oscillation probability is then given by
$$P(\gamma g)=\underset{q}{}\left|\stackrel{}{𝒢}_q|\stackrel{}{𝒱}^{(0)}+\stackrel{}{𝒱}^{(1)}\right|^2=\underset{q}{}\left|\stackrel{}{𝒢}_q|\stackrel{}{𝒱}_{\mathrm{int}}^{(1)}\right|^2,$$
(171)
that is, using (170),
$$P(\gamma g)=\underset{q}{}\left|_0^u\mathrm{\Delta }_M(u^{})\text{e}^{i\mathrm{\Delta }_m^{(q)}u^{}i_0^u^{}\mathrm{\Delta }_\lambda (u^{\prime \prime })𝑑u^{\prime \prime }}𝑑u^{}\right|^2.$$
(172)
We can check that in a homogeneous field we recover (69), i.e. the oscillation probability in the weak mixing case. Note that in the particular case of the weak mixing this method of computing the oscillation probability is shorter than the one used in the two former sections since it does not involve the determination of the eigenvalues of $``$. But, one has to assume that the probability is small compared to unity , which is not necessarily the case for instance when we are in the strong mixing regime.
### B Example of applications
As an example, we consider the mixing in a periodic magnetic field of the form $`H_0\mathrm{cos}\mathrm{\Delta }_0u`$ with $`\mathrm{\Delta }_0>0`$ for which the oscillation probability, in the weak mixing regime, is given by (172)
$$P(\gamma g)=\underset{i1}{}s_i\left|_0^u𝑑z\mathrm{\Delta }_M\mathrm{cos}(\mathrm{\Delta }_0z)\text{e}^{i\mathrm{\Delta }_m^{(r_i)}z}\text{e}^{i_0^z\mathrm{\Delta }_\lambda (v)𝑑v}\right|^2.$$
(173)
Assume that $`\mathrm{\Delta }_{\mathrm{plasma}}`$ dominates so that we can neglect the variation of $`\mathrm{\Delta }_\lambda `$ with $`z`$ then, the probability becomes
$$P(\gamma g)\mathrm{\Delta }_M^2\underset{i}{}s_i\left(\frac{1}{(\mathrm{\Delta }_m^{(r_i)}\mathrm{\Delta }_\lambda ^{()})^2}\mathrm{sin}^2\left[\frac{\mathrm{\Delta }_m^{(r_i)}\mathrm{\Delta }_\lambda ^{()}}{2}z\right]+\frac{1}{(\mathrm{\Delta }_m^{(r_i)}\mathrm{\Delta }_\lambda ^{(+)})^2}\mathrm{sin}^2\left[\frac{\mathrm{\Delta }_m^{(r_i)}\mathrm{\Delta }_\lambda ^{(+)}}{2}z\right]\right)$$
(174)
where we have kept only the resonant term, which depends on the sign of $`\mathrm{\Delta }_m^{(r_i)}\mathrm{\Delta }_\lambda `$, and where we have defined
$$\mathrm{\Delta }_\lambda ^{(\pm )}\mathrm{\Delta }_\lambda \pm \mathrm{\Delta }_0.$$
(175)
Now if $`|\mathrm{\Delta }_m^{(r_i)}\mathrm{\Delta }_\lambda ^{(+)}||\mathrm{\Delta }_M|`$ or $`|\mathrm{\Delta }_m^{(r_i)}\mathrm{\Delta }_\lambda ^{()}||\mathrm{\Delta }_M|`$ we find a strong mixing regime, meaning that because of the resonance there will exist a mode for which the probability is enhanced.
In conclusion, the important scale that fixes the photon effective mass is now $`\mathrm{\Delta }_\lambda ^{(\pm )}\mathrm{\Delta }_0`$ if we are in a regime where $`|\mathrm{\Delta }_\lambda |\mathrm{\Delta }_0`$. We can then have the same discussion as in the previous sections but with $`\beta `$ defined as
$$\beta =\frac{\mathrm{\Delta }_m}{\mathrm{\Delta }_\lambda ^{(\pm )}},$$
(176)
according to the sign of $`\mathrm{\Delta }_\lambda `$. The length scale $`\lambda _\gamma `$ is now given by $`(\omega \mathrm{\Delta }_0)^{1/2}`$ and it follows that we expect the two following effects:
1. by increasing $`\mathrm{\Delta }_0`$ we can hope to make $`\beta `$ as small as wanted and thus to get a large enhancement of the oscillation probability. What happens is that the scale $`\mathrm{\Delta }_\lambda `$ is replaced by $`\mathrm{\Delta }_0`$ and thus that a departure from the four dimensional case will be observed if $`(\omega \mathrm{\Delta }_0)^{1/2}<R`$. When the field is homogeneous, the scale $`\mathrm{\Delta }_\lambda ^1`$ is usually very large compared to $`R`$ (see § VIII and § IX), which implies that there is little hope to see any effect of the extra–dimensions. By using an inhomogeneous field, we change the scale associated with the photon effective mass wich is now gouverned by $`\mathrm{\Delta }_0^1`$ that can be tried to be lowered to a scale close to $`R`$.
2. whatever the sign of $`\mathrm{\Delta }_\lambda `$, we expect to have strong mixing occuring for all values of $`\mathrm{\Delta }_0`$ such that
$$|\mathrm{\Delta }_m^{(r_i)}(\mathrm{\Delta }_\lambda \pm \mathrm{\Delta }_0)|\mathrm{\Delta }_M.$$
By varying slowly $`\mathrm{\Delta }_0`$ or $`\omega `$, we expect to see a series of strong and weak mixing regimes.
The amplitude of these two effects will be discussed in the last section of this article.
## VIII Application to astrophysics and cosmology
Magnetic fields are observed in most astrophysical systems but the origin of galactic and cosmological magnetic fields is still unknown . A possibility is that these fields have a primordial origin since such a magnetic field can be generated in a number of early universe mechanisms such as in collisions of bubbles produced in a first order phase transition or during an inflationary phase .
The efficiency of the photon–graviton and of the photon–axion mixing depends both on the value of the magnetic field and on the spatial extension of this field, $`\mathrm{\Lambda }_c`$ say. We study the order of magnitude of these mixings on the cosmic microwave background, on pulsars and magnetars.
The required quantities for our discussion are $`\mathrm{\Delta }_M`$, $`\mathrm{\Delta }_m`$, $`\mathrm{\Delta }_{\mathrm{plasma}}`$ and $`\mathrm{\Delta }_{\mathrm{QED}}`$ respectively given by equations (29) and (35). It is usefull to rewrite these quantities numerically as
$`{\displaystyle \frac{\mathrm{\Delta }_M}{1\mathrm{cm}^1}}`$ $`=`$ $`4\times 10^{25}\left({\displaystyle \frac{H_0}{1\mathrm{G}}}\right)\text{(graviton)},`$ (177)
$`{\displaystyle \frac{\mathrm{\Delta }_M}{1\mathrm{cm}^1}}`$ $`=`$ $`2\times 10^{16}\left({\displaystyle \frac{H_0}{1\mathrm{G}}}\right)\left({\displaystyle \frac{f_{\mathrm{PQ}}}{10^{10}\mathrm{GeV}}}\right)^1\text{(axion)},`$ (178)
$`{\displaystyle \frac{\mathrm{\Delta }_m}{1\mathrm{cm}^1}}`$ $`=`$ $`{\displaystyle \frac{2.5\times 10^{28}}{\left(2.5\times 10^{15}\right)^{4/n}}}\left({\displaystyle \frac{M_D}{1\mathrm{TeV}}}\right)^{2+4/n}\left({\displaystyle \frac{\omega }{1\mathrm{eV}}}\right)^1,`$ (179)
$`{\displaystyle \frac{\mathrm{\Delta }_{\mathrm{plasma}}}{1\mathrm{cm}^1}}`$ $`=`$ $`3.6\times 10^{17}\left({\displaystyle \frac{\omega }{1\mathrm{eV}}}\right)^1\left({\displaystyle \frac{n_e}{1\mathrm{cm}^3}}\right),`$ (180)
$`{\displaystyle \frac{\mathrm{\Delta }_{\mathrm{QED}}}{1\mathrm{cm}^1}}`$ $`=`$ $`1.33\times 10^{27}\left({\displaystyle \frac{\omega }{1\mathrm{eV}}}\right)\left({\displaystyle \frac{H_0}{1\mathrm{G}}}\right)^2,`$ (181)
where we have used the facts that $`1\mathrm{eV}5\times 10^4\mathrm{cm}^1`$, $`1\mathrm{G}1.95\times 10^2\mathrm{eV}^2`$ in the natural Lorentz-Heaviside units where $`\alpha =e^2/4\pi =1/137`$ and the expression of the extra-dimensions radius
$$R=\left(2.5\times 10^{15}\right)^{2/n}10^{12}\left(\frac{M_D}{1\mathrm{Tev}}\right)^{12/n}\mathrm{eV}^1.$$
(182)
We now restrict to the case $`n=2`$.
### A Cosmic microwave background
It has been shown that the isotropy of the cosmic microwave background (CMB) puts a limit on the present value of a spatially homogeneous magnetic field to $`(B_0/1\mathrm{G})6.8\times 10^9(\mathrm{\Omega }_0h^2)^{1/2}`$ . A comparable bound has also been obtained for spatially inhomogeneous magnetic fields . We study the magnitude of the photon–graviton conversion on the two following examples:
* Large scales: we assume that we have a homogeneous magnetic field on the scale of the Hubble radius with
$$H_06\times 10^9\mathrm{G}.$$
(183)
The CMB photons are observed as a black body with a temperature of 2.7 K so that we approximatively have photons of energy
$$\omega 10^510^3\mathrm{eV}.$$
(184)
The caracteristic size of the system is the size of the Hubble radius
$$\mathrm{\Lambda }_c=3000h^1\mathrm{Mpc}10^{28}\mathrm{cm}$$
(185)
where $`h`$ is the reduced Hubble parameter. We also estimate the electronic density today to be about (see e.g. )
$$n_e10^7\mathrm{cm}^3.$$
(186)
* Degree scales: we assume a homogeneous magnetic field on the size of the Hubble radius at the last scattering surface. Since the magnetic field scales like (scale factor)<sup>2</sup> and the energy of the photon as (scale factor)<sup>-1</sup>, we assume a magnetic field of
$$H_06\times 10^3\mathrm{G}$$
(187)
and consider photons of energy
$$\omega 10^21\mathrm{eV}$$
(188)
at a redshift of $`z1000`$. The characteric size of the system is given by the Hubble radius at decoupling, i.e.
$$\mathrm{\Lambda }_c=3\times 10^{23}h^1\mathrm{cm},$$
(189)
and the electronic density at the time of decoupling is of order (see e.g. )
$$n_e10^3\mathrm{cm}^3.$$
(190)
The main idea is that, since photons are converted into either gravitons or axions, some anisotropies must be induced on the scale of homogeneity of the magnetic field, mainly because of the angular dependence of the conversion rate. The effect between a direction parallel and direction perpendicular to the magnetic field must not exceed the observed CMB temperature anisotropy. The anisotropy of the CMB temperature between the directions perpendicular to the magnetic field (where the effect of mixing is maximum) and parallel to it (where there is no mixing effect) is then of order
$$\frac{\mathrm{\Delta }T}{T}\frac{\mathrm{\Delta }T}{T}|_{}\frac{\mathrm{\Delta }T}{T}|_{||}P(\gamma g).$$
(191)
Observationally, we have the constraint that
$$\frac{\mathrm{\Delta }T}{T}<10^5.$$
(192)
From figure 3, we deduce that in both cases, $`|\mathrm{\Delta }_{\mathrm{QED}}||\mathrm{\Delta }_{\mathrm{plasma}}|`$ so that $`\mathrm{\Delta }_\lambda \mathrm{\Delta }_{\mathrm{plasma}}`$ and thus $`\beta >0`$. In the two considered regimes we have:
1. Large angular scale:
$`\alpha `$ $``$ $`6.6\times [10^{15},10^{13}]`$ (193)
$`\beta `$ $``$ $`1.1\times 10^{21}\left({\displaystyle \frac{M_D}{1\mathrm{TeV}}}\right)^4.`$ (194)
Thus, we are always in a regime where $`\alpha 1`$, $`\beta >0`$ and $`|\beta |1`$ thus we expect at most effects of order $`\alpha ^2`$ which are completely unobservable.
2. Small angular scale:
$`\alpha `$ $``$ $`6.6\times [10^{10},10^8]`$ (195)
$`\beta `$ $``$ $`1.1\times 10^{17}\left({\displaystyle \frac{M_D}{1\mathrm{TeV}}}\right)^4.`$ (196)
We are always in a regime where $`\alpha ^21`$, $`\beta >0`$ and $`|\beta |1`$ and, as in the previous case, there will be no observable effect.
From this results, with see that $`\beta `$ is always too large to have any enhancement of the probability. Moreover in both cases the oscillation length, $`\mathrm{}_{\mathrm{osc}}`$, with the lightest KK mode (as well as the oscillation length with any massive KK mode) is much smaller than $`\mathrm{\Lambda }_c`$, and the mixing angle with the graviton zero mode is very small. The effects are the same as in a standard four dimensional spacetime and thus negligible . Note that, in theory, we should have included the expansion of the universe but this will not change the result drastically. A detailed study of the photon–graviton mixing in an expanding four dimensional spacetime can be found in and a discussion of the effects of the inhomogeneity of the field in .
### B Pulsars
As proposed by many authors (see e.g. ), axions could be produced in the interior of neutron stars in nucleon–nucleon collisions. This would constitute the main cooling mechanism for these stars and thus puts limit on the axion production flux and mass. Such a production of KK gravitons in higher dimensional theories also exist and can be used to put bounds on the mass scale $`M_D`$ .
As originally proposed by Morris (see also ), this axion (and now the KK gravitons) flux may be detectable by the secondary photons produced through the mixing with these particles in the neutron star magnetosphere magnetic field. These photons have a typical energy of
$$\omega 10^4\mathrm{eV},$$
(197)
i.e. of order of the average value of the neutron star interior temperature (about 50 keV). The primary photons can be well approximated by a black body spectrum with a temperature of $`T_{NS}1\mathrm{keV}`$ typical for the surface temperature of such stars.
The idea is to detect a distortion of the star spectrum due both to the secondary photons and to the oscillation of primary photons. The typical value of the magnetic field in the neutron star magnetosphere is
$$H_010^{12}\mathrm{G}$$
(198)
on a characteristic size of the system is of order of the neutron star size
$$\mathrm{\Lambda }_c10\mathrm{km}.$$
(199)
Indeed, one cannot neglect the effect of the magnetospheric plasma and we estimate its density as
$$n_e7\times 10^2\left(\frac{H_0}{1\mathrm{G}}\right)\left(\frac{P}{1\mathrm{s}}\right)^1$$
(200)
where $`P`$ is the period of the pulsar and will be assumed to be about 1 second in the following.
According to figure 4, we deduce that $`\mathrm{\Delta }_{\mathrm{QED}}|\mathrm{\Delta }_{\mathrm{plasma}}|`$ so that $`\beta <0`$ and $`\mathrm{\Delta }_\lambda \mathrm{\Delta }_{\mathrm{QED}}`$. Then, it follows
$`\alpha `$ $``$ $`3\times 10^{14}`$ (201)
$`|\beta |`$ $``$ $`3\times 10^8\left({\displaystyle \frac{M_D}{1\mathrm{TeV}}}\right)^4`$ (202)
and we are always in a regime where $`\alpha 1`$, $`\beta <0`$ and $`|\beta |1`$ and where the characteristic size of the system is far larger than the oscillation length. We expect $`M_D`$ to be of order $`1100`$ TeV, so that we will get an amplification of order $`110^7`$ but still unobservable. From (201) we see that, contrary to the microwave background, the dominant length scale of the system is $`\mathrm{\Delta }_{\mathrm{QED}}^1`$ so that
$$\frac{\lambda _\gamma }{R}2.4\times 10^{12}\left(\frac{\omega }{1\mathrm{eV}}\right)^1\left(\frac{H_0}{1\mathrm{G}}\right)^1\left(\frac{M_D}{1\mathrm{TeV}}\right)^2$$
(203)
which is smaller than unity for the typical value of magnetic field and wavelength considered here. By going to higher frequencies and higher magnetic fields we may get a larger amplification and thus a larger effect of the extra–dimensions, this is mainly the reason why we will turn to magnetars in the following paragraph.
To finish, let us note however that in very strong magnetic fields one must take into account the photon splitting which will compete with the photon–graviton mixing. We do not discuss this effect here.
### C Magnetars and Gamma–Ray Bursts
Magnetars are pulsars with superstrong magnetic field such as SGR 1806-20 where $`H_08\times 10^{14}\mathrm{G}`$, i.e. two orders of magnitude higher than for ordinary radio pulsars. This object is associated with soft gamma ray bursts of energy of order $`1\mathrm{keV}100\mathrm{k}\mathrm{e}\mathrm{V}`$. Other examples are GB790305 and IE1841-045 and such observations are supported by models where the gamma ray bursts are triggered by cracking of the neutron star crust due to the magnetic stress .
So we consider a system such that
$$H_010^{12}10^{15}\mathrm{G},\mathrm{\Lambda }_c10\mathrm{km},$$
(204)
and
$$\omega 10^210^6\mathrm{eV}$$
(205)
Assuming that the electronic density is well approximated by (200)<sup>‡‡</sup><sup>‡‡</sup>‡‡The pulsars cited above have a period ranging from 4 to 10 seconds., we deduce that we are in the regime $`|\mathrm{\Delta }_{\mathrm{plasma}}|\mathrm{\Delta }_{\mathrm{QED}}`$ as long as
$$\left(\frac{\omega }{1\mathrm{eV}}\right)^2\left(\frac{H_0}{1\mathrm{G}}\right)4\times 10^9$$
so that we can deduce that the QED contribution always dominates in such object and then that $`\beta <0`$. On figure 5, we depict the variation of $`\alpha `$ to show that we always have $`\alpha 1`$.
Now, effects of the extra–dimensions will appear when $`\lambda _\gamma /R<1`$ and this quantity varies typically from $`10^{18}`$ to $`10^{25}`$ for $`n=2`$ assuming $`M_D1\mathrm{TeV}`$ (see equation (203)). Note also that $`\omega ^1`$ is of the order of $`R`$ and that most of the effect comes from the fact that $`\mathrm{\Delta }_{\mathrm{QED}}`$ becomes large compared with $`R^1`$. Such objects may be interesting to detect the effects of the extra–dimensions but more data and a better understanding of soft gamma–ray bursts are needed before drawing any conclusions.
## IX Laboratory experiments
Many experiments searching for light particles like axions were set up (see e.g. and for a recent review). We can classify the methods in the two following categories:
* The direct methods in which a flux of axions coming from some astrophysical source (Sun, supernovae…) is tried to be converted into photons through an external magnetic field. These transitions have been used to put bounds on astrophysical axion fluxes and coupling constant . One could think to use the same kind of experiments to put constraints in the case of a mixing with a large number of KK states. Let us first discuss the case of KK gravitons. The energy flux into KK gravitons from astrophysical object cannot exceed the bound on the energy flux in axions in the usual four dimensional case since otherwise the cooling rates of these objects will be too high . Let us then assume that the efficiency of detection is maximum and that the experiment is designed to collect all the emitted particles. Since each graviton is coupled with a much lower coupling constant than the four dimensional axion coupling constant accessible to these kinds of experiments, we do not expect that such direct detection methods will be able to see any KK graviton coming from astrophysical sources. In other words, since each KK graviton is coupled at tree level only to the photon (and not to other gravitons), there is no effect of the large number of KK states in these experiments.
* The indirect methods where one tries to detect the mixing of the photon through its effect on a photon beam in a magnetic field, both on its amplitude and polarisation. Now, the photon being coupled at tree level to all KK states, one expects a departure from the usual case. The effect on the polarisation of the beam comes from the fact that for axions only the $`\times `$ component of electromagnetic wave couples to the axions. For the gravitons, both polarisations evolve according to the same equation but, due to the QED and Cotton–Mouton birefringence, $`\mathrm{\Delta }_+\mathrm{\Delta }_\times `$ which implies a phase shift between them.
In the next paragraphs, we focuse on polarisation experiments to detect the phase shift. We first compute this phase shift in a $`D`$ dimensional spactime and discuss two kinds of experiments respectively in a static and periodic magnetic field. We must emphasize here that the magnitude of the mixing with axions depends on the free parameter $`f_{\mathrm{PQ}}`$ (in contrast with the mixing with gravitons which magnitude is fixed by the value of the Planck mass), so that these experiments may be able to put constraints on bulk axion models.
### A Phase shift in $`D`$ dimensions
Let us go back to the four dimensional case for axion. Then, from (53), we deduce that, starting from an initial state $`(A(0),G(0)=0)`$, the two polarisations $`+`$ and $`\times `$ evolve respectively as, omitting a global phase $`\omega u`$,
$$A_+(u)=\text{e}^{i\mathrm{\Delta }_+u}A_+(0),A_\times (u)=\left(\text{e}^{i\mathrm{\Delta }_\times ^{}u}\mathrm{cos}^2\vartheta +\text{e}^{i\mathrm{\Delta }_g^{}u}\mathrm{sin}^2\vartheta \right)A_\times (0).$$
(206)
Now, restricting to the weak mixing regime where $`\vartheta 1`$, we can expand the $`\mathrm{\Delta }^{}`$ defined in (52) as
$$\mathrm{\Delta }_\times ^{}\mathrm{\Delta }_\times +\vartheta ^2(\mathrm{\Delta }_\times \mathrm{\Delta }_m),\mathrm{\Delta }_g^{}\mathrm{\Delta }_m\vartheta ^2(\mathrm{\Delta }_\times \mathrm{\Delta }_m).$$
(207)
Expanding (206) to second order and taking into account (207) leads to
$$A_+(u)=\text{e}^{i\mathrm{\Delta }_+u}A_+(0),A_\times (u)=\left[1i\vartheta ^2\zeta +\vartheta ^2(\text{e}^{i\zeta }1)\right]\text{e}^{i\mathrm{\Delta }_\times u}A_\times (0)$$
(208)
with $`\zeta (\mathrm{\Delta }_\times \mathrm{\Delta }_m)u`$. We deduce that the two modes evolve relative to each other as
$$\frac{A_\times (u)}{A_+(u)}=\left[1i\vartheta ^2\zeta +\vartheta ^2(\text{e}^{i\zeta }1)\right]\text{e}^{i(\mathrm{\Delta }_\times \mathrm{\Delta }_+)u}\frac{A_\times (0)}{A_+(0)}.$$
(209)
The relative phase and amplitude of $`A_\times `$ with respect to $`A_+`$ then evolve as
$$\left|\frac{A_\times }{A_+}\right|(u)\left[12\vartheta ^2\mathrm{sin}^2\left(\frac{\zeta }{2}\right)\right]\left|\frac{A_\times }{A_+}\right|(0),\varphi (u)\left[(\mathrm{\Delta }_+\mathrm{\Delta }_\times )u\vartheta ^2(\zeta \mathrm{sin}\zeta )\right]\varphi (0).$$
(210)
When we neglect the mixing effect (i.e. $`\vartheta =0`$), there is a phase shift due to the QED and Cotton–Mouton birefringence. The extra phase shift due to the fact that only one polarisation of the axion is affected by the mixing has been used to design experiments to put constraints on the axion parameters (see e.g. ).
In the case of a graviton, the two polarisations are mixed in the same way, so that the same computation leads to
$`\left|{\displaystyle \frac{A_\times }{A_+}}\right|(u)`$ $``$ $`\left(12\left[\vartheta _\times ^2\mathrm{sin}^2\left({\displaystyle \frac{\zeta _\times }{2}}\right)\vartheta _+^2\mathrm{sin}^2\left({\displaystyle \frac{\zeta _+}{2}}\right)\right]\right)\left|{\displaystyle \frac{A_\times }{A_+}}\right|(0),`$ (211)
$`\varphi (u)`$ $``$ $`\left[(\mathrm{\Delta }_+\mathrm{\Delta }_\times )u+\vartheta _+^2\left(\zeta _+\mathrm{sin}\zeta _+\right)\vartheta _\times ^2\left(\zeta _\times \mathrm{sin}\zeta _\times \right)\right]\varphi (0)`$ (212)
(we now have to keep the index $`\lambda `$ on $`\vartheta `$ and on $`\zeta `$). In four dimension, $`\mathrm{\Delta }_m=0`$ so that the phase shift depends only on the QED and Cotton–Mouton parameters. Its amplitude is proportional to $`\mathrm{\Delta }_M`$ so that it is roughtly 10 orders of magnitude lower than for axions.
Let us now compute the phase shift in a $`D`$ dimensional spacetime. We assume that we are in the weak mixing regime and apply the method of § VII. As seen on (208), we must compute the solution of (163) up to second order. Following the same lines as for the computation leading to (170), we can show that $`\stackrel{}{𝒱}_{\mathrm{int}}^{(2)}`$ is explicitely given by
$$\stackrel{}{𝒱}_{\mathrm{int}}^{(2)}(u)=\underset{q}{}_0^u𝑑y\mathrm{\Delta }_M(y)_0^y𝑑z\mathrm{\Delta }_M(z)\text{e}^{i_y^z[\mathrm{\Delta }_m^{(q)}\mathrm{\Delta }_\lambda (x)]𝑑x}A(0)\stackrel{}{𝒜},$$
(213)
From which we deduce that, starting from a pure photon state, the polarisation $`\lambda `$ of the photon evolves as
$$A_\lambda (u)=\left[1\underset{q}{}_0^u𝑑y\mathrm{\Delta }_M(y)_0^y𝑑z\mathrm{\Delta }_M(z)\text{e}^{i_y^z[\mathrm{\Delta }_m^{(q)}\mathrm{\Delta }_\lambda (x)]𝑑x}\right]A_\lambda (0).$$
(214)
In the case of a homogeneous field, we can extract from (214) the relative phase of the polarisation $`\times `$ with respect to the polarisation $`+`$ for the case of KK gravitons and bulk axions. In the latter case, only the polarisation $`\times `$ evolves according to (214) whereas the polarisation $`+`$ evolves according to (208) so that
$`\varphi (u)`$ $`=`$ $`(\mathrm{\Delta }_+\mathrm{\Delta }_\times )u\alpha ^2{\displaystyle \underset{i1}{}}{\displaystyle \frac{s_i}{(1\beta _i^{(\times )})^2}}\left[(\mathrm{\Delta }_\times \mathrm{\Delta }_m^{(r_i)})u\mathrm{sin}(\mathrm{\Delta }_\times \mathrm{\Delta }_m^{(r_i)})u\right]\varphi (0)\text{(axion)}`$ (215)
$`=`$ $`(\mathrm{\Delta }_+\mathrm{\Delta }_\times )u{\displaystyle \underset{i1}{}}s_i\{{\displaystyle \frac{\left[(\mathrm{\Delta }_\times \mathrm{\Delta }_m^{(r_i)})u\mathrm{sin}(\mathrm{\Delta }_\times \mathrm{\Delta }_m^{(r_i)})u\right]}{(1\beta _i^{(\times )})^2}}\alpha _\times ^2`$ (217)
$`{\displaystyle \frac{\left[(\mathrm{\Delta }_+\mathrm{\Delta }_m^{(r_i)})u\mathrm{sin}(\mathrm{\Delta }_+\mathrm{\Delta }_m^{(r_i)})u\right]}{(1\beta _i^{(+)})^2}}\alpha _+^2\}\text{(graviton)}`$
with the notations used before. We split this result in three parts as
$$\varphi =\varphi _{\mathrm{QED}}+\varphi _{\mathrm{CM}}+\varphi _{KK}$$
(218)
where the two first terms are the phase shifts due to vacuum polarisation and the Cotton–Mouton effect and are obtained by setting $`\alpha =0`$ in (215-217). The third term is the specific phase shift associated with the mixing of the photon with either bulk axions or Kaluza–Klein gravitons.
### B Polarisation experiments
We now discuss two kinds of experiments designed to detect the mixing induced phase shift. As typical parameters we take $`\omega 2\mathrm{eV}`$ for the laser beam and a magnetic field which might be as strong as $`H_0=10^5`$ G. With these values, we have (for $`n=2`$)
$`{\displaystyle \frac{\mathrm{\Delta }_M}{1\mathrm{cm}^1}}4\times 10^{20},\text{(graviton)}{\displaystyle \frac{\mathrm{\Delta }_M}{1\mathrm{cm}^1}}2\times 10^{11}\left({\displaystyle \frac{f_{\mathrm{PQ}}}{10^{10}\mathrm{GeV}}}\right)^1,\text{(axion)}`$ (219)
$`{\displaystyle \frac{\mathrm{\Delta }_{\mathrm{QED}}}{1\mathrm{cm}^1}}3\times 10^{17},{\displaystyle \frac{\mathrm{\Delta }_{\mathrm{plasma}}}{1\mathrm{cm}^1}}1.3\times 10^{17}\left({\displaystyle \frac{n_e}{1\mathrm{cm}^3}}\right),{\displaystyle \frac{\mathrm{\Delta }_m}{1\mathrm{cm}^1}}1.5\times 10^4\left({\displaystyle \frac{M_D}{1\mathrm{Tev}}}\right)^4.`$ (220)
Then, the QED and plasma effects are of the same order of magnitude but, since $`\varphi _{\mathrm{QED}}\varphi _{\mathrm{CM}}`$, we neglect the Cotton-Mouton effect and define the phase shift ratio as
$$_{KK}\frac{\varphi _{KK}}{\varphi _{\mathrm{QED}}}.$$
(221)
One hopes to be able to measure a $`_{KK}`$ of order 0.1 .
#### 1 Multiple path experiments
The idea is to make a laser beam reflect between two mirrors distant of $`l`$. Since the mirror are transparent to the axions and gravitons, the phase shift after $`N`$ paths will be $`N\varphi (l)`$ so that, in the case of axions, $`_{KK}`$ is given by
$$_{KK}=\left(\frac{\mathrm{\Delta }_M}{\mathrm{\Delta }_+\mathrm{\Delta }_\times }\right)\underset{i1}{}s_i\frac{\mathrm{\Delta }_M}{\mathrm{\Delta }_\times \mathrm{\Delta }_m^{(r_i)}}\left[1\frac{\mathrm{sin}(\mathrm{\Delta }_\times \mathrm{\Delta }_m^{(r_i)})l}{(\mathrm{\Delta }_\times \mathrm{\Delta }_m^{(r_i)})l}\right].$$
(222)
The decrease of the amplitude of the photon beams due to the creation of axions is
$$\delta I4\alpha ^2N\underset{i1}{}\frac{s_i}{(1\beta _i)^2}\mathrm{sin}^2\left[\frac{\mathrm{\Delta }_\times \mathrm{\Delta }_m^{(r_i)}}{2}l\right].$$
(223)
Then, an enhancement of both the phase shift and the variation of the beam are expected due to the sums over all states. With the experimental values specified above, we have
$$\frac{\lambda _\gamma }{R}10^8\left(\frac{M_D}{1\mathrm{TeV}}\right)^2,$$
(224)
which is larger than unity. Then, we think that experiments in homogeneous magnetic fields will not probe the extra–dimensions since it will require to work with very high magnetic fields. Note that when we span the electromagnetic spectrum from the infrared to the X–ray, $`\lambda _\gamma /R`$ varies in the range
$$\frac{\lambda _\gamma }{R}4\times (10^810^2)\left(\frac{M_D}{1\mathrm{TeV}}\right)^2.$$
(225)
#### 2 Effect of a periodic field
As seen in § VII B, one can hope to enhance the mixing effect by using a periodic magnetic field. For that purpose we need the pulsation $`\mathrm{\Delta }_0`$ to dominate over $`\mathrm{\Delta }_\lambda `$ and $`\beta `$ to be small compared to unity. With the previous numerical values, the first condition rewrites as $`\mathrm{\Delta }_0>10^{17}\mathrm{cm}^1`$ and will be satisfied easily. Using (182), the second condition gives for a six dimensional spacetime
$$\mathrm{\Delta }_0^1<2.5\times 10^2\left(\frac{M_D}{1\mathrm{TeV}}\right)^4\left(\frac{\omega }{1\mathrm{eV}}\right)\mathrm{cm}.$$
(226)
As stressed before, we have a departure from the four dimensional behaviour only if $`\lambda _\gamma <R`$. Now, since the magnetic field varies on a scale $`\mathrm{\Delta }_0^1`$, $`\lambda _\gamma `$ is gouverned by $`\mathrm{\Delta }_0`$ instead of $`\mathrm{\Delta }_\lambda `$. Then an effect will appear only if we manage to create a field that can vary on scales of the order of the centimeter.
In § VII B we also quoted the possibility of having a series of strong mixing regimes, which is specific of the existence of extra–dimensions. This requires that $`|\mathrm{\Delta }_m^{(r_i)}(\mathrm{\Delta }_\lambda \pm \mathrm{\Delta }_0)|\mathrm{\Delta }_M\sqrt{s_i}`$ and can be performed either by varying $`\mathrm{\Delta }_0`$ or $`\omega `$. Let us assume that $`\mathrm{\Delta }_0`$ is fixed, since $`\mathrm{\Delta }_0\mathrm{\Delta }_M`$ we have a strong mixing regime for the pulsations defined by
$$\omega ^{(\stackrel{}{p})}=\frac{\stackrel{}{p}^2}{2R^2\mathrm{\Delta }_0}$$
(227)
with a width of
$$\delta \omega ^{(\stackrel{}{p})}=\frac{\stackrel{}{p}^2}{2R^2\mathrm{\Delta }_0}\frac{\mathrm{\Delta }_M}{\mathrm{\Delta }_0},$$
(228)
that is
$$\frac{\omega ^{(\stackrel{}{p})}}{1\mathrm{eV}}=2\times 10^4\stackrel{}{p}^2\left(\frac{M_D}{1\mathrm{TeV}}\right)^2\frac{\mathrm{\Delta }_0^1}{R},\frac{\delta \omega ^{(\stackrel{}{p})}}{1\mathrm{eV}}=10^3\stackrel{}{p}^2\left(\frac{\mathrm{\Delta }_0^1}{R}\right)^2\left(\frac{\mathrm{\Delta }_M}{1\mathrm{cm}^1}\right).$$
(229)
For instance if we assume that $`R`$ is of order of the millimeter and that we consider a field varying on the order of the meter, we get for axions that
$$\frac{\omega ^{(\stackrel{}{p})}}{1\mathrm{eV}}=2\times 10^2\stackrel{}{p}^2\left(\frac{M_D}{1\mathrm{TeV}}\right)^2,\frac{\delta \omega ^{(\stackrel{}{p})}}{1\mathrm{eV}}=2\times 10^{10}\stackrel{}{p}^2.$$
(230)
## X Conclusion
We have used the property that, in an external magnetic field, the KK gravitons and axions couple at tree level to photons to show that there exists a mixing between the photon and these particles. We have computed this mixing and compared our result to the photon-graviton and photon-axion mixings in four dimensions. The main difference comes from the fact that the mixing matrix is now infinite and that a photon couples to a large number of massive particles. We then have discussed the physical implications of these phenomena in a general $`D`$ dimensional universe. This leads us to conclude that for most astrophysical objects the effect of photon-KK gravitons mixing will be unobservable.
The main points of this study are:
* We describe how to deal with the mixing between the photon and a large number of light particles. This extends the former results to the case of $`D`$ dimensional universes where a photon can mix with all the Kaluza–Klein gravitons and possibly with bulk axions. In (6365), we gave the exact expression for the oscillation probability and we then discussed its amplitudes first qualitatively and then in a five and in a six dimensional spacetime.
* When $`\lambda _\gamma <R`$, i.e. if the caracteristic length scale associated with the photon wavelength and effective mass is smaller than the radius of the extra–dimensions, there is a departure from the four dimensional effect. Otherwise, the first KK mode is too heavy to be excited and everything, in general, reduces to the four dimensional situation.
* Two limiting regimes have been found:
+ A large radius regime where the two following behaviours can appear
- In the weak mixing regime, we have shown that the oscillation probability can be obtained by summing over the individual oscillation probabilities and is then enhanced by a factor of order $`\beta ^{n/2}`$ in comparison with the standard four dimensional case. In that case the solutions can be found either by solving the eigenvalues equation or by considering the equation of evolution as a Schrödinger equation and solving it iteratively in the interaction picture. The latter method generalises to the case of inhomogeneous magnetic fields but only applies if the oscillation probability is small compared to unity.
- In the strong mixing regime the photon mixes preferentially with a given KK modes, which is possible if the plasma effects dominate over the vacuum polarisation. In that case a complete transition is possible. We note that this effect is more likely to happen in a $`D`$ dimensional context than in a four dimensional spacetime and point out a specific effect of the extra–dimensions on the oscillation length.
+ A small radius regime where $`R\lambda _\gamma `$ so that the spacing between to KK modes is very large compared to the photon characteristic length. In that case, the probability is generally dominated by the contribution of the first state corresponding to the lightest particle and we are back to the four dimensional case. A consequence of this is that we can have an observable signature of the extra–dimensions only if $`\lambda _\gamma <R`$. In most of the systems this cannot be achieved, mainly because $`\mathrm{\Delta }_\lambda `$ is very small; the only favorable situation happens in strong magnetic fields such as in pulsars and magnetars and when we deal with a magnetic field which varies on a small enough typical scale.
* We have shown that, in the case of graviton, the effect of this mixing although enhanced is too small to be observed on the cosmic microwave background and on astrophysical objects such as pulsars. However, we point out that the effects can be larger with bulk axions.
* We discussed laboratory experiments designed for the search of axions in the light of this new framework and we computed the phase shift (215217) between the two polarisations of a photon entering a magnetic field. As for the probability, we show that the phase shift is enhanced by the existence of the KK modes. In a periodic field, we show that the effect of the extra–dimensions can be important and that there exists a series of strong and weak mixing regimes that may be observed. More work to derive bounds on $`f_{\mathrm{PQ}}`$ is however needed before drawing any conclusion.
* On a more technical level, we have shown how to diagonalise a general matrix of mixing. This kind of matrices appears in different situations in extra-dimensions physics and these results can be used in many problems and in particular for neutrino oscillations.
* To finish, let us stress some important comments. First, we have assumed the validity of the Euler–Heisenberg Lagrangian to describe the vacuum polarisation. One has also to be aware that in strong magnetic fields such as in pulsar magnetosphere, photon splitting will be in competition with the photon graviton oscillation. We have not compared the strength of these two effects but we expect the latter to be dominant in high magnetic fields, such as in the magnetosphere of magnetars. We have also concentrated on gravitational waves even if there is also production of scalar waves. They are thought to be negligible, at least in the four dimensional case . Concerning the axions, the effects may be more important than for gravitons, depending on the value of the coupling $`f_{\mathrm{PQ}}`$, but we did not reconsider the bounds on the axion parameters. Further work is needed in that direction. We also stress that in general one needs a precise determination of their mass spectrum to compute the coupling of each mode to the photon. Most of these results were obtained for $`n=1`$ or $`n=2`$ extra–dimensions for which the results respectively do not, or only weakly, depend on the cut–off of the theory. We also stress that depending on the exact physical situation the number of KK modes with which the photon can oscillate coherently can be drastically limited in comparison with the number of accessible KK modes. These decoherence effect implies that the UV cut–off of the theory $`M_{\mathrm{max}}`$ is expected to be much higher than the physical cut–off. A consequence of this is that the results obtained in the cases $`n=1`$ and $`n=2`$ can be extended to $`n>2`$ without depending on $`M_{\mathrm{max}}`$.
## Acknowledgments
We wish to thank l’École de Physique Théorique des Houches where this work was initiated, P. Binétruy, E. Dudas, J.F. Glicenstein, R. Lehoucq, M. Lemoine, J. Mourad, O. Pene, P. Peter and J. Rich for discussions.
## A Diagonalisation of $``$ in the general case
The goal of this appendix is to compute the eigenvalues and eigenvector of the matrix $``$, to diagonalise it and to explain how to compute the probability of conversion of a photon into a graviton.
We consider the $`(N+2)\times (N+2)`$ matrix defined by
$$=\left(\begin{array}{cccccc}\mathrm{\Delta }_\lambda & \mathrm{\Delta }_M& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{\Delta }_M\\ \mathrm{\Delta }_M& \mathrm{\Delta }_m^{(0)}& 0& \mathrm{}& \mathrm{}& 0\\ \mathrm{}& 0& \mathrm{}& 0& \mathrm{}& 0\\ \mathrm{}& \mathrm{}& 0& \mathrm{\Delta }_m^{(q)}& 0& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& 0\\ \mathrm{\Delta }_M& 0& \mathrm{}& \mathrm{}& 0& \mathrm{\Delta }_m^{(N)}\end{array}\right).$$
(A1)
We restrict to a finite matrix since there is a cut–off in the theory as discuseed in § IV B and all the notations are defined in § II.
Let us stress that the diagonalisation of this matrice is a purely technical point that appears often in extra-dimension physics (see e.g. ).
### 1 Characteristic polynomial
These notations being fixed, we compute the characteristic polynomial $`𝒫(x)`$ of the matrix $``$ defined by
$$𝒫(x)det\left(xI_{N+2}\right),$$
(A2)
where $`I_{N+2}`$ is the $`(N+2)\times (N+2)`$ identity matrix. Developping (A2) with respect to its first column leads to
$$𝒫(x)=\left(\mathrm{\Delta }_\lambda x\right)\underset{q=0}{\overset{N}{}}\left(\mathrm{\Delta }_m^{(q)}x\right)+\underset{q=0}{\overset{N}{}}(1)^{q+1}\mathrm{\Delta }_MD_q(x),$$
(A3)
where $`D_q`$ is the determinant of the comatrix of the element $`(q+2,1)`$ given by
$$D_q(x)=(1)^q\mathrm{\Delta }_M\underset{\mathrm{}=0,\mathrm{}q}{\overset{N}{}}\left(\mathrm{\Delta }_m^{(\mathrm{})}x\right).$$
(A4)
From (A3) and (A4), we deduce that the characteristic polynomial is given by
$$𝒫(x)=\underset{q=0}{\overset{N}{}}\left(\mathrm{\Delta }_m^{(q)}x\right)\left[\mathrm{\Delta }_\lambda x\mathrm{\Delta }_M^2\underset{q=0}{\overset{N}{}}\frac{1}{\left(\mathrm{\Delta }_m^{(q)}x\right)}\right].$$
(A5)
### 2 Eigenvalues
The characteristic eigenvalue equation $`𝒫(x)=0`$ has $`N+2`$ real solutions since $``$, being a symetric matrix, is diagonalisable. To find all his solutions, we rewrite (A5) as
$$𝒫(x)=A(x)B(x),$$
(A6)
with
$`A(x)`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{N_D}{}}}\left(\mathrm{\Delta }_m^{(r_i)}x\right)^{s_i1}`$ (A7)
$`B(x)`$ $`=`$ $`\left(\mathrm{\Delta }_\lambda x\right){\displaystyle \underset{i=1}{\overset{N_D}{}}}\left(\mathrm{\Delta }_m^{(r_i)}x\right)\mathrm{\Delta }_M^2{\displaystyle \underset{i=1}{\overset{N_D}{}}}{\displaystyle \underset{\mathrm{}=1,\mathrm{}i}{\overset{N_D}{}}}s_{\mathrm{}}\left(\mathrm{\Delta }_m^{(r_{\mathrm{}})}x\right).`$ (A8)
We have two kind of eigenvalues:
* $`\mathrm{\Delta }_m^{(r_i)}`$: they are solutions of $`A(x)=0`$ and are of order $`s_i1`$ but are not solutions of $`B(x)=0`$ so that they are solutions of $`𝒫(x)=0`$ with order $`s_i1`$ and thus are eigenvalues of $``$ of the same order.
This gives us $`_{i=1}^{N_D}(s_i1)=N+1N_D`$ eigenvalues of $``$.
* $`x_i`$: they are solutions of $`B(x)=0`$ and since $`B(x)`$ is a polynomial of order $`N_D+1`$ and since $``$ is diagonalisable, we must have $`N_D+1`$ such eigenvalues. To find them, we rewrite $`B(x)`$ as
$$B(x)=\left[\underset{i=1}{\overset{N_D}{}}\left(\mathrm{\Delta }_m^{(r_i)}x\right)\right]\left[\mathrm{\Delta }_\lambda x\mathrm{\Delta }_M^2\underset{j=1}{\overset{N_D}{}}\frac{s_j}{\left(\mathrm{\Delta }_m^{(r_i)}x\right)}\right].$$
(A9)
Let us stress that $`x_i\mathrm{\Delta }_m^{(r_i)}`$ since otherwise the cancellation occuring in the first factor is offset by a divergence in the second factor. It follows that the $`x_i`$ are solutions of
$$\mathrm{\Delta }_\lambda x=\mathrm{\Delta }_M^2\underset{i=1}{\overset{N_D}{}}\frac{s_i}{\left(\mathrm{\Delta }_m^{(r_i)}x\right)}.$$
(A10)
This solution can be found numerically but we can find the main properties of these eigenvalues graphically from which we deduce that (A10) has $`N_D+1`$ distinct solutions that we order as
$$\left(x_i\right)_{1iN_D+1}x_1<\mathrm{}<x_{n+1}.$$
(A11)
In conclusion, we have found the $`N+2`$ eigenvalues of $``$ which split in $`n`$ eigenvalues $`\mathrm{\Delta }_m^{(r_i)}`$ each with multiplicity $`s_i1`$ and in $`N_D+1`$ distinct eigenvalues $`x_i`$.
### 3 Eigenvectors
To determine the eigenvectors $`V`$ solution of
$$V=xV,$$
(A12)
we set
$$V(v,u_0,\mathrm{},u_N)$$
(A13)
so that (A12) reduces to the system
$`\mathrm{\Delta }_\lambda v+\mathrm{\Delta }_M\left({\displaystyle \underset{q=0}{\overset{N}{}}}u_q\right)=xv`$ (A14)
$`\mathrm{\Delta }_Mv+\mathrm{\Delta }_m^{(q)}u_q=xu_q.`$ (A15)
* If $`x=\mathrm{\Delta }_m^{(r_i)}`$, the eigenvectors generate a subspace of dimension $`s_i1`$ a basis of which is given explicitely by
$$V_{r_i+p}=\frac{1}{\sqrt{(p+1)(p+2)}}\left[\underset{\mathrm{}=0}{\overset{p}{}}G_{r_i+\mathrm{}}+(p+1)G_{r_i+p+1}\right],0ps_i2,$$
(A16)
where $`\{A,(G_q)_{0qN}\}`$ is the initial orthonormal basis where we have written $``$ in (A1). One can check that this is an orthonotmal family, i.e. that
$$V_{r_i+p}V_{r_j+q}=\delta _{pq}\delta _{ij}.$$
* If $`x=x_i`$, for each eigenvalue we have a subspace of dimension 1 generated by the unit vector
$$V_{x_i}=\frac{1}{\sqrt{\mathrm{\Delta }_M^2+_{q=0}^N\left(x_i\mathrm{\Delta }_m^{(q)}\right)^2}}(\frac{1}{\mathrm{\Delta }_M},\frac{1}{x_i\mathrm{\Delta }_m^{(0)}},\mathrm{},\frac{1}{x_i\mathrm{\Delta }_m^{(N)}})$$
(A17)
in the basis $`\{A,(G_q)_{0qN}\}`$. It is easy to show that they satisfy
$$V_{x_i}V_{x_j}=\delta _{ij},V_{x_i}V_{r_j+p}=0.$$
We have given the explicit form of the $`N+2`$ eigenvectors of $``$. It is worthwile noting that the eigenstates $`V_{r_i+p}`$ mix the different KK modes together while letting the photon unaffected whereas the eigenstates $`V_{x_i}`$ mix the photon with the $`N+1`$ KK gravitons.
## B Probability of oscillation in the general case
To compute the oscillation probability between a photon and gravitons in a constant magnetic field, we follow the method by Raffelt and Stodolsky and solve the equation of evolution (27) as
$$\stackrel{}{𝒱}(u)=\text{e}^{iu}\text{e}^{i\omega u}\stackrel{}{𝒱}(0),$$
(B1)
where $`\stackrel{}{𝒱}\{A,G^{(0)},\mathrm{},G^{(N)}\}`$. We decompose this vector on the eigenvectors basis as
$$\stackrel{}{𝒱}(0)=\underset{i=1}{\overset{N_D}{}}\underset{p=0}{\overset{s_i2}{}}h_{i,p}(0)V_{r_i+p}+\underset{i=1}{\overset{n+1}{}}f_i(0)V_{x_i},$$
(B2)
where $`h_{i,p}(0)`$ and $`f_i(0)`$ are $`N+2`$ coefficients. Injecting this decomposition in (B1), we obtain
$$\stackrel{}{𝒱}(z)=\left[\underset{i=1}{\overset{N_D}{}}\underset{p=0}{\overset{s_i2}{}}h_{i,p}(0)\text{e}^{i\mathrm{\Delta }_m^{(r_i)}u}V_{r_i+p}+\underset{i=1}{\overset{N_D+1}{}}f_i(0)\text{e}^{ix_iu}V_{x_i}\right]\text{e}^{i\omega u}.$$
(B3)
The probability of a photon to be converted in KK gravitons is obtained by considering the initial state $`\stackrel{}{𝒱}(0)=\{1,0,\mathrm{},0\}`$ and by computing
$$P(\gamma \gamma )=\left|\underset{q=0}{\overset{N}{}}G^{(q)}(z)\stackrel{}{𝒱}(0)\right|^2.$$
(B4)
Since only the modes associated with the eigenvalues $`x_i`$ mix with the photon, we deduce that
$$P(\gamma g)=1P(\gamma \gamma )=1\left|\underset{i=1}{\overset{N_D+1}{}}f_{x_i}^2\text{e}^{ix_iu}\right|^2,$$
(B5)
where the coefficients $`f_{x_i}`$ are
$$f_{x_i}=\left[1+\underset{q=0}{\overset{N}{}}\frac{\mathrm{\Delta }_M^2}{\left(x_i\mathrm{\Delta }_m^{(q)}\right)^2}\right]^{1/2}.$$
(B6)
## C Upper bound on $`|_J|`$ and $`|𝒢_J|`$
The goal of this appendix is to give an upper bound on the absolute value of the two functions $`_J(y)`$ and $`𝒢_J(y)`$ (defined in (120) and (C16), (C17)) when $`n=2`$. These majorations are used in section VI to determine the solution of the eigenvalues equation (119) as well as the oscillation probability (63) in a small coupling limit. We also give an upper bound on $`s_i`$. We first give a bound on
$`_{i_1,i_2}(y)`$ $``$ $`\alpha ^2{\displaystyle \underset{i=i_1}{\overset{i_2}{}}}{\displaystyle \frac{s_i}{y\beta _i}},\text{ and on}`$ (C1)
$`𝒢_{i_1,i_2}(y)`$ $``$ $`\alpha ^2{\displaystyle \underset{i=i_1}{\overset{i_2}{}}}{\displaystyle \frac{s_i}{\left(y\beta _i\right)^2}}.`$ (C2)
We recall that the $`\beta _i`$ are defined by $`\beta _i\stackrel{}{p_i}^2\beta p_i^2\beta ,`$ where $`\stackrel{}{p_i}`$ is a pair $`(n_i,m_i)`$ of integers, and $`p_i`$ is defined by the second equality. We assume all along this discussion that $`\beta `$ is positive and we order the $`\beta _i`$ as $`\beta _i<\beta _{i+1}`$. We stress that $`\beta _1=0`$.
### 1 $`0y<\beta _{i_1}<\beta _{i_2}`$
For $`ii_1`$, each $`\stackrel{}{p}_i^2`$ belongs to a unique interval
$$(p_{i_1}+k)^2p_i<(p_{i_1}+k+1)^2,$$
(C3)
where $`k`$ is an integer. $`_k`$, the number of such $`\stackrel{}{p}_i`$, is bounded by four times the surface defined by
$$p_{i_1}+k\sqrt{x_1^2+x_2^2}p_{i_1}+k+1$$
(C4)
in the real plan $`(x_1,x_2)`$, since we have at most four pairs for each unit square cell. Thus, one has
$$_k8\pi (p_{i_1}+k+1).$$
(C5)
One can then easily obtain a majoration in term of $`(p_{i_1}+k)`$:
$$_kq(p_{i_1}+k),$$
(C6)
with $`q=16\pi `$. On the other hand, for each $`p_i`$ satisfying (C3), $`|1/(y\beta _i)|`$ is lower than $`|1/(y(p_{i_1}+k)^2\beta )|`$, from which we get the upper bound on $`_{i_1,i_2}(y)`$
$$\left|_{i_1,i_2}(y)\right|\frac{\alpha ^2q}{\beta }\underset{k=0}{\overset{k_{\mathrm{max}}}{}}\frac{\sqrt{\beta }(p_{i1}+k)}{\left(\sqrt{\beta }(p_{i1}+k)\right)^2y}\sqrt{\beta }.$$
(C7)
$`k_{\mathrm{max}}`$ is defined such that $`p_{i_1}+k_{\mathrm{max}}<p_{i_2}<p_{i_1}+k_{\mathrm{max}}+1`$. The r.h.s. of (C7) is nothing else but a Riemann sum associated with the function $`f(x)=x/(x^2y)`$. Since $`f(x)`$ is decreasing for all $`x^2>y`$, we have
$`\left|_{i_1,i_2}(y)\right|{\displaystyle \frac{\alpha ^2q}{\beta }}{\displaystyle _{\sqrt{\beta }(p_{i_1}1)}^{\sqrt{\beta }(p_{i_1}+k_{\mathrm{max}})}}{\displaystyle \frac{xdx}{x^2y}}{\displaystyle \frac{\alpha ^2q}{2\beta }}\mathrm{ln}\left({\displaystyle \frac{\beta (p_{i_1}+k_{\mathrm{max}})^2y}{\beta (p_{i_1}1)^2y}}\right).`$ (C8)
Let us emphasize that (C8) assumes implicetly that $`y<\beta (p_{i_1}1)^2`$. Otherwise, the contribution of the $`\beta _i`$ such that $`\beta p_{i_1}^2\beta _i<\beta (p_{i_1}+1)^2`$ in the sum (C1) can be bounded by
$$\frac{\alpha ^2q}{\beta p_{i_1}^2y}p_{i_1}.$$
(C9)
Since $`\beta (p_{i_1}1)^2y<\beta p_{i_1}^2`$, we deduce that
$$p_{i_1}\left(\frac{\sqrt{y}}{\sqrt{\beta }}+1\right),$$
(C10)
so that we get the majoration
$`|_{i_1,i_2}(y)|{\displaystyle \frac{\alpha ^2q}{(\beta p_{i_1}^2y)}}\left({\displaystyle \frac{\sqrt{y}}{\sqrt{\beta }}}+1\right)+{\displaystyle \frac{\alpha ^2q}{\beta }}\mathrm{ln}\left({\displaystyle \frac{\beta (p_{i_1}+k_{\mathrm{max}})^2y}{\beta p_{i_1}^2y}}\right).`$ (C11)
With similar arguments, on can show that, when $`y<\beta (p_{i_1}1)^2`$, $`|𝒢_{i_1,i_2}(y)|`$ is bounded by
$`|𝒢_{i_1,i_2}(y)|`$ $``$ $`{\displaystyle \frac{\alpha ^2q}{\beta }}{\displaystyle _{\sqrt{\beta }(p_{i_1}1)}^{\sqrt{\beta }(p_{i_1}+k_{\mathrm{max}})}}{\displaystyle \frac{xdx}{(x^2y)^2}}{\displaystyle \frac{\alpha ^2q}{2\beta }}\left({\displaystyle \frac{1}{\beta (p_{i_1}1)^2y}}{\displaystyle \frac{1}{\beta (p_{i_1}+k_{\mathrm{max}})^2y}}\right),`$ (C12)
and that otherwise
$`|𝒢_{i_1,i_2}(y)|`$ $``$ $`{\displaystyle \frac{\alpha ^2q}{(\beta p_{i_1}^2y)^2}}\left({\displaystyle \frac{\sqrt{y}}{\sqrt{\beta }}}+1\right)+{\displaystyle \frac{\alpha ^2q}{2\beta }}\left({\displaystyle \frac{1}{\beta p_{i_1}^2y}}{\displaystyle \frac{1}{\beta (p_{i_1}+k_{\mathrm{max}})^2y}}\right).`$ (C13)
The bounds (C11) and (C13) are valid also in the case where $`y<\beta (p_{i_1}1)^2`$.
### 2 $`0\beta _{i_1}<\beta _{i_2}<y`$
In that case, following the same line of reasoning, we obtain respectively for $`_{i_1,i_2}`$ and $`𝒢_{i_1,i_2}`$ and any $`y`$ satisfying the above condition
$$|_{i_1,i_2}(y)|\frac{\sqrt{y}\alpha ^2q}{\sqrt{\beta }(y\beta p_{i_2}^2)}+\frac{\alpha ^2q}{2\beta }\mathrm{ln}\left(\frac{y\beta (p_{i_2}k_{\mathrm{max}})^2}{y\beta p_{i_2}^2}\right),$$
(C14)
$$|𝒢_{i_1,i_2}(y)|\frac{\sqrt{y}\alpha ^2q}{\sqrt{\beta }(y\beta p_{i_2}^2)^2}+\frac{\alpha ^2q}{\beta }\left(\frac{1}{y\beta p_{i_2}^2}\frac{1}{y\beta (p_{i_2}k_{\mathrm{max}})^2}\right),$$
(C15)
### 3 bound on $`_J`$ and $`𝒢_J`$
We now give a bound on the expression defined by
$`_J(y)`$ $``$ $`\alpha ^2{\displaystyle \underset{i=1,iJ}{\overset{N_D}{}}}{\displaystyle \frac{s_i}{y\beta _i}},`$ (C16)
$`𝒢_J(y)`$ $``$ $`\alpha ^2{\displaystyle \underset{i=1,iJ}{\overset{N_D}{}}}{\displaystyle \frac{s_i}{\left(y\beta _i\right)^2}},`$ (C17)
For $`J`$ defined by
$$iJ,|y\beta _i||y\beta _J|.$$
(C18)
One has thus $`iJ,|y\beta _i|\beta /2`$, and in particular $`\beta p_{J+1}^2y>\beta /2`$ and $`y\beta p_{J1}^2>\beta /2`$. Using equations (C11C15) as well as the value of $`p_{\mathrm{max}}`$ given<sup>\**</sup><sup>\**</sup>\**Note that the energy cut–off $`p_{\mathrm{max}}`$ has to be regarded as a maximum one fixed by the underlying quantum regularisation of the theory. However, decoherence effects can reduce drastically the number of KK modes that have to be taken into account. Thus, the bounds on $`𝒬`$ and $`𝒬^{}`$ in a more realistic case are likely to be lower than the one given here. in (33), one obtains easily the following bound on $`|_J|`$ and $`|𝒢_J|`$ <sup>\*†</sup><sup>\*†</sup>\*†Similar bounds would also apply to $`(y)`$ and $`𝒢(y)`$ when $`y`$ is so that $`i[1,N_D],|y\beta _i|\beta /2`$..
$`|_J|`$ $``$ $`𝒬{\displaystyle \frac{\alpha ^2}{\beta ^2}}\text{ sup }(\sqrt{\beta y},𝒬^{})`$ (C19)
$`|𝒢_J|`$ $``$ $`𝒬{\displaystyle \frac{\alpha ^2}{\beta ^2}}\text{ sup }({\displaystyle \frac{\sqrt{y}}{\sqrt{\beta }}},2)`$ (C20)
where $`𝒬`$ is a constant of order $`10^3`$ and $`𝒬^{}`$ a constant of order 10.
### 4 bound on $`s_i`$
In the real plane the euclidian distance between two different pairs of integers is bounded by $`1`$, so that the number of pairs of integers on a given closed curve is always lower than the length of this curve. It is then easy to obtain the following bound on $`s_i`$ which represents then number of different pairs of integers on a circle of radius $`p_i`$.
$$s_i2\pi p_i=2\pi \sqrt{\frac{\beta _i}{\beta }}.$$
(C21)
To finish, let us note that the properties of the series $`s_i`$ have been studied by Gauss around 1800, see e.g. \[http://mathworld.wolfram.com/rn.html\] for details and references. |
warning/0002/math0002182.html | ar5iv | text | # Solutions of the Einstein-Dirac Equation on Riemannian 3-Manifolds with Constant Scalar Curvature. 11footnote 1Supported by the SFB 288 of the DFG.
## 1 Introduction
Consider a Riemannian spin manifold of dimension $`n3`$ and denote by $`D`$ the Dirac operator acting on spinor fields. A solution of the Einstein-Dirac equation is a spinor field $`\psi `$ solving the equations
$$Ric\frac{1}{2}Sg=\pm \frac{1}{4}T_\psi ,D(\psi )=\lambda \psi .$$
Here $`S`$ denotes the scalar curvature of the space, $`\lambda `$ is a real constant and $`T_\psi `$ is the energy-momentum tensor of the spinor field $`\psi `$ defined by the formula
$$T_\psi (X,Y)=(X_Y\psi +Y_X\psi ,\psi ).$$
The scalar curvature $`S`$ is related to the eigenvalue $`\lambda `$ and the length of the spinor field $`\psi `$ by the formula
$$S=\pm \frac{\lambda }{n2}|\psi |^2.$$
In \[KimF\] we introduced the weak Killing equation for a spinor field $`\psi ^{}`$:
$$_X\psi ^{}=\frac{n}{2(n1)}dS(X)\psi ^{}+\frac{2\lambda }{(n2)S}Ric(X)\psi ^{}\frac{\lambda }{n2}X\psi ^{}+\frac{1}{2(n1)S}XdS\psi ^{}$$
Any weak Killing spinor $`\psi ^{}`$ (WK-spinor) yields a solution $`\psi `$ of the Einstein-Dirac equation after normalization
$$\psi =\sqrt{\frac{(n2)|S|}{|\lambda ||\psi ^{}|^2}}\psi ^{}.$$
In fact, in dimension $`n=3`$ the Einstein-Dirac equation is essentially equivalent to the weak Killing equation (see \[KimF\]). Up to now the following 3-dimensional Riemannian manifolds admitting WK-spinors are known:
1. the flat torus $`T^3`$ with a parallel spinor;
2. the sphere $`S^3`$ with a Killing spinor;
3. two non-Einstein Sasakian metrics on the sphere $`S^3`$ admitting WK-spinors. The scalar curvature of these two left-invariant metrics equals $`S=1\pm \sqrt{5}`$.
The aim of this paper is to classify all Riemannian 3-manifolds with constant scalar curvature and admitting a solution of the Einstein-Dirac equation. In particular, we will prove the existence of a one-parameter family of left-invariant metrics on $`S^3`$ with WK-spinors. This family contains the two non-Einstein Sasakian metrics with WK-spinors on $`S^3`$, but does not contain the standard sphere $`S^3`$ with Killing spinors. Moreover, any simply-connected, complete Riemannian manifold $`N^3S^3`$ with WK-spinors and constant scalar curvature is isometric to a space of this one-parameter family. In order to formulate the result precisely, we fix real parameters $`K,L,MR`$ and denote by $`N^3(K,L,M)`$ the 3-dimensional, simply-connected and oriented Riemannian manifold defined by the following structure equations:
$$\omega _{12}=K\sigma ^3,\omega _{13}=L\sigma ^2,\omega _{23}=M\sigma ^1,$$
or, equivalently:
$$d\sigma ^1=(LK)\sigma ^2\sigma ^3,d\sigma ^2=(M+K)\sigma ^1\sigma ^3,d\sigma ^3=(LM)\sigma ^1\sigma ^2.$$
The 1-forms $`\sigma ^1,\sigma ^2,\sigma ^3`$ are the dual forms of an orthonormal frame of vector fields. Using this frame the Ricci tensor of $`N^3(K,L,M)`$ is given by the matrix
$$Ric=\left(\begin{array}{ccc}2KL& 0& 0\\ 0& 2KM& 0\\ 0& 0& 2LM\end{array}\right).$$
Main Theorem: Let $`N^3S^3`$ be a complete, simply-connected Riemannian manifold with constant scalar curvature $`S0`$. If $`N^3`$ admits a WK-spinor, then $`N^3`$ is isometric to $`N^3(K,L,M)`$ and the parameters are a solution of the equation
$$K^2L(LM)^2M+L^3M^3+KL^2M^2(ML)+K^3(LM)(L+M)^2=0()$$
Conversely, any space $`N^3(K,L,M)`$ such that $`(K,L,M)(0,0,0)`$ is a solution of $`()`$ admits two WK-spinors for one and only one WK-number $`\lambda `$. With respect to the fixed orientation of $`N^3(K,L,M)`$ we have the two cases:
$$\lambda =+\frac{S}{2\sqrt{2}}\sqrt{\frac{S}{S^22|Ric|^2}}\text{if }K<M$$
$$\lambda =\frac{S}{2\sqrt{2}}\sqrt{\frac{S}{S^22|Ric|^2}}\text{if }M<K.$$
The spaces $`N^3(K,L,M)`$ are isometric to $`S^3`$ equipped with a left-invariant metric.
Remark: If the parameters $`K=M`$ coincide, the solution of the equation $`()`$ is given by
$$L=\frac{1}{4}K(1\sqrt{5}),L=\frac{1}{4}K(1+\sqrt{5})$$
and we obtain the Ricci tensors
$$Ric=\left(\begin{array}{ccc}\frac{1}{2}K^2(\sqrt{5}1)& 0& 0\\ 0& 2K^2& 0\\ 0& 0& \frac{1}{2}K^2(\sqrt{5}1)\end{array}\right)$$
or
$$Ric=\left(\begin{array}{ccc}\frac{1}{2}K^2(1+\sqrt{5})& 0& 0\\ 0& 2K^2& 0\\ 0& 0& \frac{1}{2}K^2(1+\sqrt{5})\end{array}\right).$$
The non-Einstein-Sasakian metrics on $`S^3`$ occur for the parameter $`K=1`$ (see \[KimF\]).
Remark: Using the standard basis of the Lie algebra $`\text{so}(3)`$ we can write the left-invariant metric of the space $`N^3(K,L,M)`$ in the following way:
$$\left(\begin{array}{ccc}\frac{1}{|ML||K+M|}& 0& 0\\ & & \\ 0& \frac{1}{|KL||ML|}& 0\\ & & \\ 0& 0& \frac{1}{|KL||K+M|}\end{array}\right).$$
The equation $`()`$ is a homogeneous equation of order six. The transformation $`(K,L,M)(\mu K,\mu L,\mu M)`$ corresponds to a homothety of the metric. Therefore - up to a homothety - the moduli space of solutions is a subset of the real projective space $`P^2(R)`$ given by the equation $`()`$. This subset is a configuration of six curves in $`P^2(R)`$ connecting the three points $`[K:L:M]=[1:0:0],[0:1:0],[0:0:1]`$ corresponding to flat metrics.
In particular, we have constructed two paths of solutions of the Einstein-Dirac equation deforming the non-Einstein Sasakian metrics on $`S^3`$.
## 2 The integrability condition for the Einstein-Dirac equation in dimension $`n=3`$.
The spinor bundle of a 3-dimensional Riemannian manifold is a complex vector bundle of dimension two. Moreover, there exists a quaternionic structure commuting with the Clifford multiplication by real vectors (see \[F\]). Consequently, in case of a real WK-number $`\lambda `$, the corresponding space of WK-spinors is a quaternionic vector space. In the spinor bundle let us introduce the metric connection $`^\lambda `$ given by the formula
$$_X^\lambda \psi :=_X\psi \frac{3}{4}dS(X)\psi \lambda \left\{\frac{2}{S}Ric(X)X\right\}\psi \frac{1}{4S}XdS\psi $$
and denote by $`\mathrm{\Omega }^\lambda `$ its curvature form. Then we obtain the following
Proposition 1: Let $`N^3`$ be a simply-connected 3-dimensional Riemannian manifold and suppose that the scalar curvature $`S0`$ does not vanish. Then the following conditions are equivalent:
1. $`N^3`$ is a non-trivial solution of the Einstein-Dirac equation with real eigenvalue $`\lambda `$;
2. $`N^3`$ admits a WK-spinor with real WK-number $`\lambda `$;
3. $`N^3`$ admits two WK-spinors with real WK-number $`\lambda `$;
4. The curvature $`\mathrm{\Omega }^\lambda 0`$ vanishes identically.
If the scalar curvature $`S0`$ is constant, the condition $`\mathrm{\Omega }^\lambda 0`$ has been investigated and yields algebraic equations involving the Ricci tensor and its covariant derivative (see \[KimF\], Theorem 8.3.). In order to formulate the integrability condition, we denote by $`X\times Y`$ the vector cross product of two vectors $`X,YT(N^3)`$. For brevity, let us introduce the endomorphism $`T:T(N^3)T(N^3)`$ given by the formula
$$T(X)=\underset{i=1}{\overset{3}{}}e_i\times (_{e_i}Ric)(X),$$
which will be used in the proof of the main Theorem.
Theorem 1 (see \[KimF\]): Let $`N^3`$ be a simply-connected 3-dimensional Riemannian manifold with constant scalar curvature $`S0`$. $`N^3`$ admits a solution of the Einstein-Dirac equation with real eigenvalue $`\lambda `$ if and only if the following three conditions are satisfied:
1. $`8\lambda ^2\{S^22|Ric|^2\}=S^3;`$
2. $`8\lambda ^2\{SRic(X)2RicRic(X)\}4\lambda ST(X)S^2Ric(X)=0;`$
3. $`8\lambda ^2\{2Ric(X)SX\}\times \{2Ric(Y)SY\}+8\lambda S\{(_XRic)(Y)(_YRic)(X)\}`$
$`+S^3X\times Y=2S^2\underset{i<j}{}\{R_{jY}\delta _{iX}+R_{iX}\delta _{jY}\}e_i\times e_j.`$
## 3 Proof of the Main Theorem
We fix an orthonormal frame $`e_1,e_2,e_3`$ of vector fields on $`N^3`$ consisting of eigenvectors of the Ricci tensor:
$$Ric=\left(\begin{array}{ccc}A& 0& 0\\ 0& B& 0\\ 0& 0& C\end{array}\right).$$
Denote by $`\sigma ^1,\sigma ^2,\sigma ^3`$ the dual frame and consider the connection forms $`\omega _{ij}=e_i,e_j`$ of the Levi-Civita connection. The structure equations of the Riemannian manifold $`N^3`$ are
$$d\omega _{12}=\omega _{13}\omega _{32}+\frac{CAB}{2}\sigma ^1\sigma ^2$$
$$d\omega _{13}=\omega _{12}\omega _{23}+\frac{BAC}{2}\sigma ^1\sigma ^3$$
$$d\omega _{23}=\omega _{21}\omega _{13}+\frac{ABC}{2}\sigma ^2\sigma ^3$$
and the covariant derivative $`Ric`$ is given by the matrix of 1-forms
$$Ric=\left(\begin{array}{ccc}dA& (AB)\omega _{12}& (AC)\omega _{13}\\ (AB)\omega _{12}& dB& (BC)\omega _{23}\\ (AC)\omega _{13}& (BC)\omega _{23}& dC\end{array}\right).$$
Using the third equation of Theorem 1 we obtain
$$(_{e_i}Ric)(e_j)(_{e_j}Ric)(e_i),e_i=0$$
and, consequently,
$$dA(e_1)=dA(e_2)=dB(e_1)=dB(e_3)=dC(e_2)=dC(e_3)=0.$$
Since $`A+B+C=S`$ is constant, we conclude that any eigenvalue $`A,B,C`$ of the Ricci tensor is constant, too. The second equation of Theorem 1 yields the condition that all elements outside the diagonal of the (1,1)-tensor $`T`$ are zero:
$$(AB)\omega _{12}(e_1)=0=(AB)\omega _{12}(e_2)$$
$$(CA)\omega _{13}(e_1)=0=(CA)\omega _{13}(e_3)$$
$$(BC)\omega _{23}(e_2)=0=(BC)\omega _{23}(e_3).$$
First, we discuss the generic case that $`A,B,C`$ are pairwise different. Then there exist numbers $`K,L,M`$ such that
$$\omega _{12}=K\sigma ^3,\omega _{13}=L\sigma ^2,\omega _{23}=M\sigma ^1.$$
The parameter triples $`\{A,B,C\}`$ and $`\{K,L,M\}`$ are related via the structure equations by the formulas
$$A=2KL,B=2KM,C=2LM.$$
The first and second equation of Theorem 1 become equivalent to the following system of algebraic equations:
* $`\lambda =\pm {\displaystyle \frac{S}{2\sqrt{2}}}\sqrt{{\displaystyle \frac{S}{S^22|Ric|^2}}}`$ ;
* $`2S(S^22|Ric|^2)\{(AC)L+(BA)K\}^2=S(SA2A^2)A(S^22|Ric|^2)`$
$`2S(S^22|Ric|^2)\{(CB)M+(AB)K\}^2`$ $`=`$ $`S(SB2B^2)B(S^22|Ric|^2)`$
$`2S(S^22|Ric|^2)\{(BC)M+(CA)L\}^2`$ $`=`$ $`S(SC2C^2)C(S^22|Ric|^2).`$
We solve this system of algebraic equations with respect to the parameters $`K,L,M`$. It turns out that the equations 2’. can be written in the form
$$P_i(K,L,M)Q(K,L,M)=0,$$
$`(1i3)`$, where the polynomials $`P_1,P_2,P_3`$ and $`Q`$ are given by the formulas
$`P_1(K,L,M)`$ $`=`$ $`(KL^2+L^2M+K^2(L+M))^2`$
$`P_2(K,L,M)`$ $`=`$ $`(KM^2+LM^2+K^2(L+M))^2`$
$`P_3(K,L,M)`$ $`=`$ $`(LM(L+M)+K(L^2+M^2))^2`$
$`Q(K,L,M)`$ $`=`$ $`K^2L(LM)^2M+L^3M^3+KL^2M^2(ML)+K^3(LM)(L+M)^2.`$
The real solutions of $`P_1=P_2=P_3=0`$ are the pairs $`\{K=0,L=0\}`$ (the flat metric) and $`\{K=M,L=M\}`$ (the space of positive constant curvature). Therefore, we proved that a 3-dimensional complete, simply-connected manifold $`N^3`$ with constant scalar curvature $`S0`$ and different eigenvalues of the Ricci tensor is isometric to one of the spaces $`N^3(K,L,M)`$, where the parameters $`K,L,M`$ are solutions of the equation $`Q(K,L,M)=0`$. These spaces satisfy the conditions 1. and 2. of Theorem 1 and, moreover, a simple computation yields the result that condition 3. of Theorem 1 is satisfied, too. We next discuss the case that two of the eigenvalues $`A,B,C`$ coincide, for example, $`A=CB`$. Then we obtain again
$$\omega _{12}=K\sigma ^3,\omega _{23}=M\sigma ^1,$$
but there is no condition for the connection form $`\omega _{13}`$. We compute the matrix of the $`(1,1)`$-tensor $`T`$:
$$T=\left(\begin{array}{ccc}(BC)K& 0& 0\\ 0& (CB)(K+M)& 0\\ 0& 0& (BC)M\end{array}\right).$$
Since the scalar curvature $`S`$ as well as the eigenvalues $`A=C,B`$ of the Ricci tensor are constant, the second equation of Theorem 1 yields that $`K`$ and $`M`$ are constant and, moreover, coincide:
$$K=M=\text{const.}$$
In case of $`K=M=0`$ we have $`\omega _{12}=\omega _{23}=0`$ and $`A=C`$. In particular, the Ricci tensor is parallel, $`Ric=0`$. Therefore, in this case $`N^3`$ is a Ricci-parallel 3-dimensional manifold admitting a WK-spinor. Then $`N^3`$ is either flat or a space of constant positive curvature (see \[KimF\], Theorem 8.2.). Finally, we consider that the case of $`K=M=1`$, i.e., $`\omega _{12}=\sigma ^3`$ and $`\omega _{23}=\sigma ^1`$. Differentiating the equation $`\omega _{12}=\sigma ^3`$, we obtain
$$\omega _{13}\omega _{32}\frac{B}{2}\sigma ^1\sigma ^2=d\omega _{12}=d\omega ^3=\omega _{31}\sigma ^1+\omega _{32}\sigma ^2$$
$$\frac{B}{2}\sigma ^1\sigma ^2=\sigma ^1\sigma ^2.$$
Consequently, $`B=2`$ and the tensors $`T`$ and $`Ric`$ are given by the matrices
$$T=\left(\begin{array}{ccc}2C& 0& 0\\ 0& 2(C2)& 0\\ 0& 0& 2C\end{array}\right),Ric=\left(\begin{array}{ccc}C& 0& 0\\ 0& 2& 0\\ 0& 0& C\end{array}\right).$$
The second condition of Theorem 1 yields the equations $`(S=2+2C)`$:
$`8\lambda ^2(SC2C^2)4\lambda S(2C)S^2C`$ $`=`$ $`0`$
$`8\lambda ^2(2S8)+8\lambda S(2C)2S^2`$ $`=`$ $`0`$
Solving these equations with respect to $`\lambda `$ and $`C`$ we obtain the three solutions:
1. $`C=2`$ and $`\lambda =\pm \frac{3}{2}`$. Then $`N^3`$ is isometric to $`S^3`$.
2. $`C=1`$ and $`\lambda =0`$. Then the scalar curvature $`S=0`$ is zero.
3. $`C=\frac{1}{2}(1\pm \sqrt{5})`$ and $`\lambda =1\pm \frac{\sqrt{5}}{2}`$. These metrics are the non-Einstein Sasakian metrics on $`S^3`$ admitting WK-spinors (see \[KimF\]). The corresponding space is contained in the family $`N^3(K,L,M)`$.
We have discussed all possibilities and, therefore, we have finished the proof of the main Theorem.
## 4 The moduli space of solutions
The moduli space of all 3-dimensional Riemannian manifolds with constant scalar curvature $`S0`$ and WK-spinors is given by the triples $`\{K,L,M\}`$ of real numbers satisfying the equation of order six $`Q(K,L,M)=0`$. The polynomial $`Q`$ is symmetric in $`\{K,L,M\}`$. Denote by
$$\gamma _1=KL+M,\gamma _2=KL+KMLM,\gamma _3=KLM$$
the elementary symmetric functions of these variables. Then we have
$$Q=4\gamma _1\gamma _2\gamma _3\gamma _2^34\gamma _3^2.$$
Consider the projective variety $`V_CP^2(C)`$ defined by the homogeneous polynomial $`Q`$:
$$V_C=\{[K:L:M]P^2(C):Q(K,L,M)=0\}.$$
$`V_C`$ has three singular points:
$$V_C^{\mathrm{sing}}=\{[1:0:0],[0:1:0],[0:0:1]\}$$
and these points correspond to the flat metric. We will now parametrize the variety $`V_C`$ by two meromorphic functions defined on a smooth Riemann surface. $`V_C`$ is given by the equation $`(K=1)`$:
$$Q(1,L,M)=L^3(M1)^2(M+1)+L^2M(1+M)^2LM^2(1+M)M^3=0.$$
Let us introduce the variables
$$a=MLLM,b=(LM)LM.$$
Then we obtain $`Q(1,L,M)=a^3+4b(1+a)`$ and the equation defining the variety $`V_C`$ becomes much simpler:
$$b=\frac{1}{4}\frac{a^3}{1+a}.$$
Next we consider a square root of $`a+1`$ and we solve the equations
$$z^21=a=MLLM,\frac{1}{4}\frac{(z^21)^3}{z^2}=b=(LM)LM$$
with respect to $`L`$ and $`M`$. Then we obtain four solution pairs $`\{L,M\}`$ depending on the variable $`z`$. For example,
$$L(z)=\frac{(1+z)(12z+z^2+\sqrt{(1+z)(1+3z5z^2+z^3}))}{4z}$$
$$M(z)=\frac{(1+z)(12z+z^2+\sqrt{(1+z)(1+3z5z^2+z^3}))}{4z}.$$
The polynomial
$$(z+1)(1+3z5z^2+z^3)=(z+1)(z1)(z+(2+\sqrt{5}))(z+(2\sqrt{5}))$$
has four different zeros. The square root $`\sqrt{(1+z)(1+3z5z^2+z^3)}`$ is a meromorphic function on the compact Riemann surface of genus $`g=1`$. Consequently, there exists a torus $`C/\mathrm{\Gamma }`$ and elliptic functions $`L,M:C/\mathrm{\Gamma }P^1(C)`$ such that the components of the variety $`V_C\backslash V_C^{\mathrm{sing}}`$ are parametrized by $`L`$ and $`M`$:
$$V_C=\{[1:L(z):M(z)]:zC/\mathrm{\Gamma }\}.$$
The functions $`LM`$ and $`LM`$ are given by the formulas:
$$LM=\frac{(1+z)(z1)^2}{2z},LM=\frac{(1+z)^2(z1)}{2z}$$
The moduli space we are interested in coincides with the real points of the projective variety $`V_C`$. If $`K=0`$, the only solutions of the equation $`Q(0,L,M)=0`$ are $`L=0`$ or $`M=0`$, i.e., the points $`[0:1:0]`$ and $`[0:0:1]`$. Therefore we can parametrize the moduli space by the parameter $`MR`$ solving the equation $`Q(1,L,M)=0`$ with respect to $`L=L(M)`$. In this way we obtain a configuration of six curves in $`P^2(R)`$ connecting the three singular points of $`V_C`$ (see the figure in the Introduction). However, we obtain geometrically different metrics on $`S^3`$ only for two curves parametrized by the real parameter $`0M\mathrm{}`$. The graphs of the function $`L_\pm (M)`$ are given in the following figure:
(Figure 1)
The functions $`L_\pm (M)`$ are monotone and tend to $`\pm 1`$ in case that $`M`$ tends to infinity. Let us discuss the geometric invariants of these metrics. The graph of the scalar curvatures $`S_\pm (M)`$ depending on $`M`$ is given by the next figure:
(Figure 2: The scalar curvatures)
Next we plot the eigenvalues $`A_\pm (M),B_\pm (M),C_\pm (M)`$ of the Ricci tensor for both families of metrics:
(Figure 3: The eigenvalues of the Ricci tensor for $`L_+(M)`$)
(Figure 4: The eigenvalues of the Ricci tensor for $`L_{}(M)`$)
In dimension $`n=3`$ the number
$$\lambda ^2(D)[\text{vol}(N^3)]^{\frac{2}{3}}$$
is a homothety invariant, where $`\lambda (D)`$ is an eigenvalue of the Dirac operator. In case of a WK-spinor we have
$$\lambda ^2=\frac{1}{8}\frac{S^3}{S^22|Ric|^2}$$
and, therefore, we obtain the formula
$$\lambda ^2\text{vol}^{\frac{2}{3}}=\frac{1}{8}(2\pi ^2)^{\frac{2}{3}}\frac{S^3}{S^22|Ric|^2}\frac{1}{\{|KL||ML||K+M|\}^{\frac{2}{3}}}.$$
The next figures contain the graph of $`\lambda ^2\text{vol}^{\frac{2}{3}}(M)`$ depending on the parameter $`M`$ for both families of metrics.
(Figure 5: $`\lambda ^2\text{vol}^{\frac{2}{3}}`$ in case of $`L_+(M)`$)
(Figure 6: $`\lambda ^2\text{vol}^{\frac{2}{3}}`$ in case of $`L_{}(M)`$)
Finally, let us discuss the behaviour of the rational function
$$\mathrm{\Psi }=\frac{L^2}{KM}$$
on the variety $`V_CP^2(C)`$. It turns out that $`\mathrm{\Psi }`$ has simple zeros at the singular points $`[1:0:0]`$ and $`[0:0:1]`$. Indeed, solving the equation defining $`V_C`$ with respect to $`L=L(M)`$ $`(K=1)`$ we obtain
$$\underset{M0}{lim}\frac{L^2(M)}{M}=0,\underset{M0}{lim}\frac{d}{dM}\left(\frac{L^2(M)}{M}\right)=1.$$
The third singular point $`[0:1:0]`$ is a pole of order two. In the regular part of $`V_C`$ the function $`\mathrm{\Psi }`$ has 12 ramification points. Among them 10 points are first order ramification points. The ramification points of order two are the points
$$[K:L:M]=[1:\frac{1}{4}(1\pm \sqrt{5}):1].$$
These parameters correspond precisely to the non-Einstein Sasakian metrics on $`S^3`$ admitting solutions of the Einstein-Dirac equation. |
warning/0002/cond-mat0002468.html | ar5iv | text | # Field-Induced Crossover and Colossal Magnetoresistance in La0.7Pb0.3MnO3
## Abstract
A field-induced crossover is observed in the resistivity ($`\rho `$) and magnetization (M) of a La<sub>0.7</sub>Pb<sub>0.3</sub>MnO<sub>3</sub> single crystal. The field-dependence of $`\rho `$ and M suggests that a small spin-canted species with mean-field-like interactions dominates at low fields (H), whereas, individual spins and 3D Ising/Heisenberg models describe the high-H behavior rather well. Around the ferromagnetic transition, an H-induced destruction of the small spin-canted magnetic polarons is accompanied by large magnetoresistance.
Colossal magnetoresistance (CMR), a large resistivity ($`\rho `$) drop induced by the magnetic field (H) in A<sub>1-x</sub>A’<sub>x</sub>MnO<sub>3</sub> (where A = rare earth elements and A’= divalent cations such as Ba, Sr, Ca or Pb), has attracted great attention recently . It is widely recognized that mean-field models with the double exchange (DE) and the Jahn-Teller (JT) distortion can not quantitatively describe CMR . Various phase-segregation models, which treat manganites as metallic ferromagnet-clusters embedded in an insulating paramagnet-matrix, were subsequently proposed . According to these models, CMR results from the formation of percolating paths between these clusters, whose concentration and size increase with H. This proposition is consistent with various observations in manganites with small (A, A’) , where the large size-mismatch between (A, A’) and Mn favors mesoscopic charge-segregation. Such a percolating description, however, may not be proper for manganites with good size-match and optimal doping (whose resistivity is metallic above T<sub>C</sub>), although experimental data in these compounds, e.g. the pair-density-function of La<sub>1-x</sub>Sr<sub>x</sub>MnO<sub>3</sub> and the non-Curie-Weiss susceptibility ($`\chi `$) in La<sub>1-x</sub>Pb<sub>x</sub>MnO<sub>3</sub> , strongly suggest the existence of local magnetic structures. The role of these structures in the CMR is still unknown. This uncertainty is the motivation behind the present investigation.
In this work, the bulk-magnetization (M) and resistivity ($`\rho `$) of a La<sub>0.7</sub>Pb<sub>0.3</sub>MnO<sub>3</sub> single-crystal are measured. An H-induced crossover is observed at a threshold field H<sub>T</sub>(T) with the following properties: 1) $`\rho `$ is independent of both H and M for H $`<`$ H<sub>T</sub>, but decreases exponentially with M for H far above H<sub>T</sub>; 2) $`\rho `$(M,T)/$``$T peaks at T<sub>C</sub> as that predicted in critical scaling only if H $`<`$ H<sub>T</sub>, but becomes zero with a universal $`\rho `$ $``$ exp(-M/M$`_\text{o}`$) at higher fields; 3) In the critical region, M(T, H) follows a mean-field-like scaling below H<sub>T</sub>, but switches to that of 3D-Heisenberg/Ising models far above; 4) The low-field Curie constant, after corrections for the critical fluctuations, is three times larger than that of isolated Mn<sup>3+</sup>/Mn<sup>4+</sup>’s. The data, therefore, suggest a robust but small magnetic-structure with canted spins, similar to the one-site small polarons observed in La<sub>1-x</sub>Sr<sub>x</sub>MnO<sub>3</sub> . We propose that the CMR in these less-distorted manganites is closely related to an H-induced destruction of the short-range spin-canted correlations.
Large single crystals of La<sub>0.7</sub>Pb<sub>0.3</sub>MnO<sub>3</sub> have been grown in sealed Pt crucibles by slowly cooling molten mixtures of La<sub>2</sub>O<sub>3</sub> (99.99%, Aldrich), MnO<sub>2</sub> (Spec. pure), PbO (Spec. pure), and PbF<sub>2</sub> (99.99%, Johnson- Mathey). A sample of 2$`\times `$2$`\times `$5 mm<sup>3</sup> was cut from a crystal and used here. The measurements of M and $`\rho `$ were made inside a Quantum Design 5-T SQUID magnetometer, and were carried out with the field increasing from 0.01 to 5 T stepwise at a fixed temperature. Care has been taken to reduce temperature fluctuations and field relaxation, and the H was calculated based on the demagnetization factor obtained from the low-field dM/dH at T<sub>C</sub>.
X-ray diffraction analysis at room temperature shows a cubic symmetry of Pm3m without noticeable orthorhombic or rhombohedral distortions. The lattice parameters are a = b = c = 3.8938(4) Å and $`\alpha `$ = $`\beta `$ = $`\gamma `$ = 90$`^\text{o}`$. Both the symmetry and the lattice parameter demonstrate an ion-size match better than that in La<sub>0.7</sub>Sr<sub>0.3</sub>MnO<sub>3</sub>, although the related lattice-disorder (due to the larger size-difference between La and Pb) seems to suppress its T<sub>C</sub> slightly . The composition was analyzed using a JEOL JXA 8600 electron microprobe. A stoichiometry La:Pb:Mn:O = 0.679(9):0.303(6):1:3.06(5) is observed over the whole surface of the crystal. The composition spread is well within our instrumental resolution of 1-3% for cations and $``$ 5% for oxygen.
The $`\rho `$(T,M) observed is shown in figure 1, with the corresponding raw data shown in the inset. Two very different regions can be clearly identified. At the high field limit, $`\rho `$ decreases with M almost universally as $`\rho `$ $``$ exp(-M/M$`_\text{o}`$) (the solid line in Fig. 1) with a fitting parameter M$`_\text{o}`$ = 23 emu/g $``$ 0.2 T. Below a threshold field H<sub>T</sub>, however, $`\rho `$ becomes M-independent. At 320 K, for example, $`\rho `$ varies less than 0.2% when M increases to half of the saturated moment M$`_{\text{sat}}`$ at 0.14 T. It is interesting to note that H<sub>T</sub> should scale with $`\epsilon ^{\beta +\gamma }`$(where $`\epsilon `$, $`\beta `$ and $`\gamma `$ are the reduced temperature $`\epsilon =`$(T-T<sub>C</sub>)/T<sub>C</sub> and two critical exponents ) in scaling models, and a 30-times increase is expected at $`\epsilon `$ = -0.04, -0.004, and 0.04 (T = 320, 332, and 346 K). However, the measured values, defined as the field where ($`\rho `$(0)-$`\rho `$(H<sub>T</sub>))/$`\rho `$(0) = 0.01, vary more slowly, i.e. 0.15, 0.04, and 0.1 T, respectively (filled triangles in Fig. 1). The crossover, therefore, is unlikely due to critical fluctuations alone. In fact, various scaling calculations show that the magnetoresistance below T<sub>C</sub> should be proportional to either M<sup>2</sup> or $`\left|\text{H}\right|`$ up to a characteristic field of k<sub>B</sub>T<sub>C</sub>/$`\mu `$ $``$ 200 T (where $`\mu `$ is the moment of the spins), which is in disagreement with our data.
Previously, similar exponential M-dependence of $`\rho `$ has been observed in manganite films . This dependence was interpreted by assuming ln$`\rho `$ $``$ -t, t $``$ cos($`\theta _{ij}`$/2) and M $``$ cos($`\theta _{ij}`$/2) in the framework of DE, where t is the transfer-integral and $`\theta _{ij}`$ is the angle between adjacent spins. In doing so, an implicit $`\theta _{ij}`$ = 2$`\overline{\theta }`$ was assumed ($`\overline{\theta }`$ is the average angle with H), and all Mn<sup>3+/4+</sup> spins were treated as independent .
The data, therefore, demonstrate that the spins around a carrier in La<sub>0.7</sub>Pb<sub>0.3</sub>MnO<sub>3</sub> are strongly correlated below H<sub>T</sub>. A field of 0.15 T at 320 K, for example, substantially aligns the spins to raise M to $``$ 65% of its saturation value M$`_{\text{sat}}`$, but leaves $`\rho `$ almost unchanged. The $`\overline{\theta }`$ = arccos(M/M$`_{\text{sat}}`$) decreased with H by more than 40$`^\text{o}`$ while the variation of $`\theta _{ij}`$ is negligible. This demonstrates a short-range spin-correlation. The spins within a mean free path ($``$ 10 Å ) from a carrier have to rotate as a whole under the field, i.e. a robust short-range spin-correlation.
It should be pointed out that the M-increase is not caused by domain rotations. The second-order transition of La<sub>0.7</sub>Pb<sub>0.3</sub>MnO<sub>3</sub>, which is expected from its lattice match and demonstrated in our scaling analysis, eliminate domains, which mark the discontinuity of order parameters. The residual ones, if any, should be isolated around defects, and involve only a small fraction of spins. Experimentally, the H-induced nucleation of domains around T<sub>C</sub> should appear as a jump-like (or S-shape) feature at low H in the Arrott plot , but could not be observed in our data (Fig.2b).
This H-independent $`\rho `$ at low-fields drops quickly with the decrease of temperature, and d$`\rho `$/dT $``$ $`\rho `$/$``$T$`_M`$ peaks at T<sub>C</sub> = 334(1) K. A scaling fit of d$`\rho `$/dT $`\left|\epsilon \right|^\alpha `$ then leads to $`\alpha `$ = 0.03(5) within $`\left|\epsilon \right|`$ $``$ 0.005. Although the uncertainties in both T<sub>C</sub> and $`\alpha `$ are large, the low-field $`\rho `$ seems to follow the same scaling law as that in Ni. At high H, however, $`\rho `$/$``$T$`_M=0`$, and the $`\rho `$ depends on M universally as reported previously . An H-induced crossover, therefore, is clear in the T-dependence of $`\rho `$ also.
The data demand a robust canted-FM spin-correlation around the carriers below H<sub>T</sub>. In fact, the resistivity drops by 50% or more when the field increases from H<sub>T</sub> to 5 T, which requires a large $`\theta _{ij}`$ (estimated $``$ 10-40$`^\text{o}`$ in the temperature range investigated) below H<sub>T</sub>. The fact that these magnetic structures are destroyed by a H $``$ H<sub>T</sub>, as shown in $`\rho `$(T,M) discussed above and the M(T,H) below, supports the conclusion. It is interesting to note that the pair-density function (PDF) data in similar La<sub>1-x</sub>Sr<sub>x</sub>MnO<sub>3</sub> show that the holes form one-site small polarons at 300-350 K, but as moving carriers at 10 K, although both temperatures are below the corresponding T<sub>C</sub> . This can happen in the DE model only if the $`\theta _{ij}`$ between the hole (Mn<sup>4+</sup>) and the adjacent Mn<sup>3+</sup> is large enough to make the hopping time longer than the time-window of PDF. A spin-canting with a T-dependent $`\theta _{ij}`$ will offer a natural interpretation for the unusual PDF data. The $`\rho `$ and M data here, therefore, present an indirect evidence for spin canting, which was not observed before . Detailed neutron diffraction investigations are needed to explore this issue.
A canted-FM correlation may occur if there is competition between various magnetic interactions, as expected in early models . This view, however, has been challenged recently . It has been pointed out that mechanisms other than those of de Gennes may be needed under certain conditions if canted states exist. Spin canting, in a sense, is incompatible with mesoscopic segregation between FM and AFM regions, and would not be expected in highly distorted manganites. Some two-orbital model calculations, however, show that phase-segregation appears in optimally doped manganites only if the JT strength ($`\lambda `$) is large enough . The almost perfect structure of La<sub>0.7</sub>Pb<sub>0.3</sub>MnO<sub>3</sub>, therefore, may eliminate mesoscopic charge segregations and stabilize the spin canting.
In order to further explore the magnetic structures, the bulk magnetization is analyzed in the critical region using the scaling function:
$$(H/M)^{1/\gamma }=a\epsilon +bM^{1/\beta }$$
(1)
where a and b are two critical amplitudes, and $`\beta `$ and $`\gamma `$ are the critical exponents defined before. It is known that the values of $`\beta `$ and $`\gamma `$ depend on the range of interactions . They will follow the mean-field theory with $`\beta `$ $``$ 0.5 and $`\gamma `$ $``$ 1 for a long range interaction, but 3D-Heisenberg or -Ising models with $`\beta `$ $``$ 0.33 and $`\gamma `$ $``$ 1.3 for a nearest-neighbor one . Previously, the critical phenomena have been investigated with rather divergent results. Sample differences were suggested as the main reason for the discrepancy . We propose that the crossover discussed above may play a major role. The magnetization data between -0.04 $``$ $`\epsilon `$ $``$ 0.04 was, therefore, analyzed separately above and below 1 T in Fig. 2 and its Inset, respectively. The field of 1 T, which is shown in the figures as filled diamonds, is chosen between H<sub>T</sub> and the field where $`\rho `$ starts to depend on M exponentially. The parameters $`\beta `$ and $`\gamma `$ in Eq. 1 were first adjusted so that the isotherms in the Arrott plots of M<sup>1/β</sup> vs. (H/M)<sup>1/γ</sup>, form a system of parallel straight lines (Fig. 2) . The intercepts Y = M<sub>S</sub>($`\epsilon `$) and X = 1/$`\chi `$$`|`$<sub>H=0</sub> were deduced, and the Kouvel-Fisher (KF) equations of Y/(dY/dT) = (T-T<sub>C</sub>)/$`\beta `$ and X/(dX/dT) = (T-T<sub>C</sub>)/$`\gamma `$ plotted against T . The intercepts of the KF plots were then used as T<sub>C</sub>, and their slopes as 1/$`\beta `$ and 1/$`\gamma `$, respectively. Finally, the scaling law M= H<sup>1/δ</sup> at T<sub>C</sub>, the constraint $`\delta `$ -1-$`\gamma `$/$`\beta `$ = 0, as well as the scaling hypothesis of M(H,T)/$`\epsilon ^\beta `$ = f<sub>±</sub>\[H$`\epsilon ^{\beta \gamma }`$\] were used to verify the scaling over the field ranges specified. These constraints are satisfied over the field range specified within the data uncertainty, which is the standard deviation from three consecutive measurements. Two sets of rather different fitting parameters are indeed obtained. T<sub>C</sub> = 334.4(1) K, $`\beta `$ = 0.33(1), and $`\gamma `$ = 1.27(2) (Fit A) were obtained for M(H $``$ 1 T) (Fig. 2). The M(H $``$ 1 T) observed, however, significantly differs from this Arrott plot (Fig. 2). While the low-field deviations below a few hundred Oe, such as those in the Inset of Fig. 2, has been observed in most ferromagnets due to some yet-to-be-found reasons , those up to 1 T are rather unusual. The M(H $``$ 1 T), on the other side, can be fit well with T<sub>C</sub> = 336.5(2) K, $`\beta `$ = 0.50(2), and $`\gamma `$ = 1.0(1) (Fit B) (Inset, Fig. 2). It is interesting to note that the critical exponents of Fit A are very close to those predicted by 3D Ising/Heisenberg models, and those of Fit B are in good agreement with the mean-field-model predictions. This seems to be consistent with the local structures proposed above. With larger physical sizes, the interaction between the magnetic polarons may be mean-field like, but the 3D-Heisenberg/Ising models will be more proper when the species are individual spins. We, therefore, associate the unusual scaling behavior with the crossover observed in $`\rho `$(T,M). It should also be noted that the Fit A is similar to the previous result in La<sub>0.7</sub>Sr<sub>0.3</sub>MnO<sub>3</sub>, where a field up to 5 T was used , and Fit B to that in La<sub>0.8</sub>Sr<sub>0.2</sub>MnO<sub>3</sub>, where the measurements were limited to below 1 T . It is possible that the canted spin-correlation and the related field-induced crossover are common features in the less distorted manganites, although their net effects may depend on the structure parameters, e.g. the tolerance factor and the lattice-disorder .
The susceptibility $`\chi `$ = M/H was measured up to 400 K at various H to estimate the size of the structure. The (T,H) phase-space, where the canted-spins exists, was first explored. The dM/dH at 340 K and 370 K was calculated, and fit with the parameters in Fit B (Inset, Fig. 3). While only the M below 1 T shows reasonable agreement with the fit at 340 K, all M fits well at 370 K. The field needed to induce the crossover appears to increase with T, and the spin-canting structure persists above T<sub>C</sub> at low-fields. The generalized Curie-Weiss equation of $`\chi `$ = C$`\epsilon ^\gamma `$/T was then used to estimate the number of the spins involved, where $`\gamma `$’ and C are an effective critical exponent and Curie constant, respectively . This equation is a high-order expansion of the 3D Heisenberg model, and fits the magnetization of Ni between 10<sup>-3</sup> $``$ $`\epsilon `$ $``$ 0.8 excellently with the fit parameters consistent with those obtained from other methods . The $`\chi _{H=0}`$ deduced from the Arrott plots of Fig 2 and its inset was used below 340 K, and the M/H directly measured above . The obtained Curie constant corresponds to an effective moment of 4.5 and 12 $`\mu _B`$ for the Fits A and B, respectively, where $`\mu _B`$ is the Bohr magneton. Although an accurate analysis is difficult due to the limited temperature range, these values suggest that the spin canting structure involves only the nearest neighbors, in agreement with that proposed above.
In conclusion, a field induced crossover is observed in La<sub>0.7</sub>Pb<sub>0.3</sub>MnO<sub>3</sub> for both $`\rho `$(T, M) and M(T, H) around T<sub>C</sub>. Above 1 T, M(T, H) varies with $`\epsilon `$ similar to the predictions of 3D-spin models, and $`\rho `$ $``$ exp(-M/M<sub>o</sub>) can be accommodated with the DE model. At lower fields, however, M(T, H) follows mean-field-like fluctuations, and $`\rho `$ is independent of M as if the spin alignment around the holes is no longer affected by H, indicative of a small but robust canted spin-structure, likely the one-site magnetic polaron previously reported. Our data suggest that the destruction of the correlation, which has been ignored in most models, is an important mechanism in the CMR of the less distorted manganites.
###### Acknowledgements.
We thank Dr Z. Y. Weng for a careful review of the manuscript, D. K. Ross and R. Bontchev for the microprobe and XRD measurements. The work is supported in part by NSF Grant No. DMR 9804325, the T. L. L. Temple Foundation, the John and Rebecca Moores Endowment and the State of Texas through TCSUH at University of Houston; and by the Director, Office of Energy Research, Office of Basic Energy Sciences, Division of Materials Sciences of the U.S. Department of Energy under Contract No. DE-AC03-76SF00098 at LBL. |
warning/0002/astro-ph0002526.html | ar5iv | text | # Absorption-Line Probes of Gas and Dust in Galactic Superwinds
## 1 Introduction
By now, it is well-established that galactic-scale outflows of gas (sometimes called ‘superwinds’) are a ubiquitous phenomenon in the most actively star- forming galaxies in the local universe (Heckman, Lehnert, & Armus 1993; Dahlem 1997; Bland-Hawthorn 1995). They are powered by the energy deposited in the interstellar medium by massive stars via supernovae and stellar winds. Over the history of the universe, outflows like these may have polluted the intergalactic medium with metals (e.g. Giroux & Shull 1997) and dust (Alton, Davis, & Bianchi 1999; Aguirre 1999a,b), heated and polluted the intracluster medium (e.g. Gibson, Loewenstein, & Mushotzky 1997; Ponman, Cannon, & Navarro 1999), and may have established the mass-metallicity relation and radial metallicity gradients in galactic spheroids (e.g. Carollo & Danziger 1994). However, the astrophysical relevance of superwinds can not be reliably assessed without first understanding their physical, dynamical, and chemical properties. To date, most of the pertinent information has come from observations of the X-ray emission produced by the hot gas (e.g. Dahlem, Weaver, & Heckman 1998) or the optical line-emission produced by the warm gas (e.g. Lehnert & Heckman 1996a).
In the present paper, we take a complementary approach, and discuss an extensive body of new data that probes the outflowing gas via its interstellar absorption-lines. This technique has some important advantages. First, since the gas is seen in absorption against the background starlight, there is no possible ambiguity as to the sign (inwards or outwards) of any radial flow that is detected. Second, the strength of the absorption will be related to the column density of the gas. In contrast, the X-ray or optical surface-brightness of the emitting gas is proportional to the emission-measure. Thus, the absorption-lines more fully probe the whole range of gas densities in the outflow, rather than being strongly weighted in favor of the densest material (which may contain relatively little mass). Finally, provided that suitably- bright background sources can be found, interstellar absorption-lines can be used to study outflows in high-redshift galaxies where the associated X-ray or optical emission may be undetectably faint. This promise has already been realized in the case of the ’Lyman Dropout’ galaxies, where the kinematic signature of outflows is clear in their rest-frame UV spectra (Franx et al 1997; Pettini et al 1998, 1999).
A few pioneering studies have already detected interstellar absorption-lines from superwinds in local starburst galaxies. Phillips (1993) discussed spatially-resolved optical spectroscopy of the $`NaI`$ “D” doublet in NGC 1808, showing that an outflow of gas at velocities of up to 700 km s<sup>-1</sup> could be traced over a region several kpc in size, coincident with a region of extra- planar dust plumes. Several recent papers (Lequeux et al. 1995; Heckman & Leitherer 1997; Sahu & Blades 1997; Kunth et al. 1998; Gonzalez-Delgado et al 1998a) have detected blueshifted interstellar absorption-lines in $`HST`$ and $`HUT`$ UV spectra of a handful of starburst galaxies, implying outflows of metal- bearing gas at velocities of 10<sup>2</sup> to 10<sup>3</sup> km s<sup>-1</sup>.
Most of the strong resonance lines of cosmically-abundant ions are found in the UV (e.g. Morton 1991; Savage & Sembach 1996), and so must be studied with $`HST`$ or $`FUSE`$ in local starbursts. In the present program, we have instead exploited the relatively greater sensitivity and availability of ground-based telescopes at visible wavelengths to study a large sample of starbursts/superwinds using the $`NaI`$ doublet at $`\lambda `$$`\lambda `$5890,5896 Å. In a few cases, we have also observed the $`KI`$ $`\lambda `$$`\lambda `$7665,7699 Å doublet, since it probes gas with nearly the same ionization state as $`NaI`$, but is more likely to be optically thin (the $`K`$ abundance is down from $`Na`$ by a factor of 15, while the two doublets have similar oscillator strengths). The ionization potentials of $`NaI`$ and $`KI`$ are only 5.1 eV and 4.3 eV respectively, so these species should primarily probe the $`HI`$ and $`H_2`$ ISM phases. Observations in vacuum-UV will be required to study the hotter and more highly ionized gas in absorption.
## 2 The Data
### 2.1 Sample Selection
Table 1 lists the salient properties of the 32 objects in our sample. These objects have been drawn from the larger samples of infrared-selected galaxies studied by Armus, Heckman, & Miley (1989 - hereafter AHM) and by Lehnert & Heckman (1995 - hereafter LH95). The specific selection criteria used in these in two programs are described in detail in these references. Briefly, AHM selected on the basis of far-IR flux and very warm far-IR color-temperatures. LH95 selected on the basis of far-IR flux, moderately warm far-IR color-temperatures, and high galaxy inclination ( disk galaxies seen within $``$30 of edge-on). The AHM and LH95 samples overlap in galaxy properties, but the former preferentially selects more powerful and more distant objects.
AHM measured the equivalent widths of the $`NaD`$ lines in their sample using low-resolution spectra. The galaxies from AHM were selected for the present program on the basis of the brightness of their nucleus at $``$ 5900Å and the equivalent width of $`NaD`$. The galaxies from LH95 had no prior measures of the $`NaD`$ line, and were selected based primarily on their proximity and availability at the time of the observations.
Of the 32 objects, only 3 are classified as bona fide AGN of the basis of their optical spectra: the type 2 Seyferts NGC7582 and Mrk273 and the highly peculiar AGN IRAS11119+3257. No published classification exists for IRAS10502-1843. The remaining objects are optically classified as $`HIInuclei`$ or $`LINER^{}s`$, and are presumed to be primarily powered by dusty starbursts (Lutz, Veilleux, & Genzel 1999). Even in the Seyfert galaxies NGC7582 and Mrk273, optical spectra show the presence of a young stellar population in the nucleus (Schmitt, Storchi- Bergmann, & Cid-Fernadez 1999; Gonalzez-Delgado, Heckman, & Leitherer 2000). Thus, with the possible exception of IRAS11119+3257 and IRAS10502-1843, the objects in our sample all contain powerful starbursts that can drive superwinds.
### 2.2 Observations
The observations were undertaken during the period from 1988 through 1994 using three different facilities: the 4-meter Blanco Telescope with the Cassegrain Spectrograph at $`CTIO`$, the 4-meter Mayall Telescope with the RC Spectrograph at $`KPNO`$, and the 2.5-meter Dupont Telescope with the Modular Spectrograph at the Las Campanas Observatory. Various spectrograph configurations were used at each observatory, the details of which are listed in Table 2. The spectral resolution used to study the $`NaD`$ lines ranged from 55 to 170 km s<sup>-1</sup>. While low by the standards of interstellar absorption-line studies, the resolution was good enough to cleanly resolve the $`NaD`$ lines in most cases (deconvolved line widths of 100 to 600 km s<sup>-1</sup> \- see below).
### 2.3 Data Reduction & Analysis
The spectra were all processed using the standard LONGSLIT package in $`IRAF`$. All the data were bias-subtracted using the overscan region of the chip and then flat-fielded using observations of either a quartz-lamp- illuminated screen inside the dome or of a quartz lamp inside the spectrograph. The spectra were then rectified using observations of bright stars to determine and remove the distortion perpendicular to the dispersion direction and observations of a $`HeNeAr`$ arc lamp to determine the two- dimensional dispersion solution. The zero-points in the wavelength scale were verified by measuring the wavelengths of strong night-sky emission- lines. Corrections to the heliocentric reference-frame were computed for the spectra. The spectra were then sky-subtracted by interactively fitting a low-order polynomial along the spatial direction, column-by-column. For a few of the galaxies we obtained relatively low-dispersion spectra that will be used to measure the reddening (Balmer decrement and continuum color). These data were flux-calibrated using observations of spectrophotometric standard stars, and otherwise were reduced in the same way as the other data.
The spectra were analyzed using the interactive SPLOT spectral fitting package in $`IRAF`$. In all cases, a one-dimensional ‘nuclear’ spectrum was extracted, covering a region with a size set by the slit width and summed over 3 to 5 pixels in the spatial direction (typically 2 by 3 arcsec). The corresponding linear size of the projected aperture is generally a few hundred to a few thousand parsecs in these galaxies (median diameter 700 pc). This is a reasonable match to the typical sizes of powerful starbursts like these (e.g. Meurer et al 1997; Lehnert & Heckman 1996b). Prior to further analysis, each 1-D spectrum was normalized to unit intensity by fitting it with, and then dividing it by, a low-order polynomial. These nuclear spectra are shown in Figure 1. Similar one-dimensional spectra for off- nuclear regions were extracted over the spatial region with adequate signal- to-noise in the continuum for each galaxy. The primary focus of the present paper is on the nuclear spectra, but we will describe the results obtained in the off-nuclear bins when these are particularly illuminating or interesting.
Given the relatively low resolution of our spectra, the likelihood that the observed $`NaD`$ line profile contains many unresolved and/or blended kinematic sub-components, and the saturated nature of the $`NaD`$ lines, we have chosen to parameterize the lines as simply as possible. Thus, for each extracted spectrum, we have fit the $`NaD`$ doublet with a single pair of Gaussians, constrained to have the same line width and a wavelength separation appropriate to the redshifted doublet. In a few objects, the adjacent $`HeI\lambda `$5876 nebular emission-line was strong and broad enough to slightly contaminate the blue half of the $`NaD\lambda `$5890 profile. For these cases, we first fit and subtracted the $`HeI`$ emission line. Only the parameters of the stronger member of $`KI`$ doublet at 7665 Å were measured. The weaker member at 7699 Å was detected, with a strength consistent with the doublet being optically-thin (i.e. an equivalent width ratio of $``$2:1).
We have not attempted a rigorous determination of the measurement uncertainties associated with these data. The relatively high signal-to-noise in the nuclear spectra (typically better than 30:1 per pixel) means that the uncertainties in the measured quantities will be dominated in most cases by systematic effects due to the contamination of the $`NaD`$ line by weaker stellar photospheric features (whose ubiquity likewise makes it difficult to determine the true continuum level to use in the line-fitting) and by the mismatch in profile shape between the actual data and the single Gaussian component used to fit each member of the doublet (see section 3.2 below). The most straightforward way to estimate the measurement uncertainties is to compare the measurements for the 11 galaxies in the sample for which we have more than one independent spectrum (taken at a different position angle). We have done so, and the results are reported in the Notes to Table 3.
## 3 Results
### 3.1 The Stellar vs. Interstellar Contribution
Before using the $`NaD`$ line to diagnose conditions in the starburst galaxies, it is imperative to establish that the line is primarily interstellar in origin in these galaxies. The $`NaD`$ line is strong in the spectra of cool stars, reaching a peak strength in the range from K3 through M0 (see Jacoby, Hunter, & Christian 1984). These stellar types can make a significant contribution to the optical spectrum of a starburst galaxy. First, the oldest underlying population in the galaxy bulge will have a dominant contribution from K -type giants (indeed the $`NaD`$ line is one of the strongest stellar absorption-lines in optical spectra of early-type galaxies and bulges - e.g. Heckman 1980; Bica et al 1991). Second, for starbursts with ages greater than about 10 Myr, cool supergiants make a significant contribution to the optical and near-IR light (Bruzual & Charlot 1993; Leitherer et al 1999).
We have therefore tried to estimate empirically what fraction of the measured equivalent width of the $`NaD`$ doublet is contributed by late-type stars in our sample galaxies. To do so, we have considered other absorption-lines that are conspicuous in the spectra of late-type stars and galactic nuclei, but which arise from highly-excited states and are therefore of purely stellar origin (i.e. they are not resonance lines like $`NaD`$). The best-studied example is the $`MgI`$ b-band at 5174Å. The strength of this line is well-correlated with the strength of the $`NaD`$ line in spectra of the nuclei of normal galaxies ( Bica et al 1991; Heckman 1980) and in stars (Jacoby, Hunter, & Christian 1984). We have used the latter two data sets to determine a best-fit to the correlation: $`W_{NaD}`$ $``$ 0.75 $`W_{Mgb}`$. The measured strength of $`Mgb`$ in our galaxies (from AHM, Veilleux et al 1995, or our own unpublished spectra) was then used to predict the equivalent width of the stellar contribution to the observed $`NaD`$ line. We have also compared our data to spectra of K giant stars obtained during the same observing runs listed in Table 2. Rather than measuring the strengths of a few particular stellar features, we have used the entire ensemble of features in the range between about 5750 and 6450 Å to estimate by-eye the fractional stellar contribution to the $`NaD`$ line. The agreement between these two methods is generally satisfactory (the predicted $`NaD`$ stellar equivalent widths agree on-average to $``$ 0.1 dex). Heckman & Lehnert (2000) have measured the fraction of the red continuum contributed by cool stars for the seven nuclei in the present sample having the highest quality detection of the interstellar component of the $`NaD`$ line. They find that this fraction is 20 to 30%, consistent with the rough estimates reported here.
As listed in Table 3, the estimated stellar contribution to the observed $`NaD`$ line in our sample galaxies ranges from negligible ($`<`$10%) to substantial ($`>`$70%), with hints of a bimodal distribution. Thus, rather than attempting a very uncertain direct correction for the effects of the stellar contribution, we have taken the simpler approach of dividing our sample into two bins: the strong-stellar-contamination objects (‘SSC’) in which stars produce $``$40% of the measured $`NaD`$ equivalent-width, and the interstellar-dominated objects (‘ISD’) in which the stellar contribution is $``$ 30%. In the discussion to follow, we will see that the $`NaD`$ lines in the two sub-samples have significantly different properties, which can be readily understood as reflecting the relative importance of the stellar and interstellar components.
### 3.2 Kinematics
The most robust indicator of an outflow is the presence of interstellar absorption-lines that are significantly blueshifted with respect the the systemic velocity of the galaxy ($`v_{sys}`$). Thus, we have first compiled the best available measures of $`v_{sys}`$ for our galaxies The velocities of the nuclear emission-lines are potentially affected by radial gas flows and are not always reliable indicators of $`v_{sys}`$ (Mirabel & Sanders 1988; Lehnert & Heckman 1996a). We have therefore determined $`v_{sys}`$ from (in order of preference) spatially-resolved galactic rotation curves, global mm-wave $`CO`$ line profiles, nuclear stellar velocities, global $`HI\lambda `$21cm emission-line profiles, and optical nuclear emission- line velocities (only used for 4 objects). See Table 1 for details and the estimated uncertainties.
The results are shown in Figure 2. For the $`ISD`$ subsample there is strong trend for the centroid of the $`NaD`$ feature to be blueshifted with respect to $`v_{sys}`$. Specifically, 11 of the 18 $`ISD`$ nuclei have $`NaD`$ blueshifts $`\mathrm{\Delta }`$$`v`$ greater than 100 km s<sup>-1</sup> (hereafter the ‘outflow sources’). In addition, while the nuclear $`NaD`$ absorption-line in NGC1808 lies close to $`v_{sys}`$, the galaxy exhibits strongly blueshifted absorption over a several-kpc-scale region along its minor axis (Phillips 1993). We therefore include it as a 12th member of the outflow sample. The net blueshifts in the $`ISD`$ nuclei are in the range $`\mathrm{\Delta }`$$`v`$ $``$ 100 to 300 km s<sup>-1</sup>, (with the exception of IRAS11119+3257). <sup>1</sup><sup>1</sup>1IRAS11119+3257 has perhaps the most peculiar optical spectrum of any ultra-luminous system. It shows very broad (1500 km s<sup>-1</sup>) Balmer, \[OIII\]$`\lambda `$$`\lambda `$4959,5007, FeII, HeI, and \[OI\]$`\lambda `$6300 emission-lines. It appears to be a member of the “I Zw 1” class of quasars (e.g. Phillips 1976), or possibly related to Mrk 231. It is very compact (barely resolved) in optical images (Armus, Heckman, & Miley 1987). The $`NaD`$ absorption profile is complex, with a strong narrow system that is blueshifted by 934 km s<sup>-1</sup>, and a weaker system blueshifted by 1410 km s<sup>-1</sup>. See Table 3 and Figure 1. In contrast to these large blueshifts, no net redshifts greater than 100 km s<sup>-1</sup> are observed in the $`ISD`$ sample. Moreover, none of the 14 members of the $`SSC`$ sample show a net $`NaD`$ blueshift or redshift that is greater than 70 km s<sup>-1</sup>. This is consistent with expectations that the velocity of the nuclear stellar $`NaD`$ component will be very close to $`v_{sys}`$.
The $`NaD`$ linewidths in the $`ISD`$ and $`SSC`$ subsamples are also significantly different (Figure 3). The lines are relatively narrow in the $`SSC`$ subsample ($`W`$ 100 to 300 km s<sup>-1</sup>, with a median of 180 km s<sup>-1</sup>), and much broader in the $`ISD`$ nuclei ($`W`$ $``$ 150 to 600 km s<sup>-1</sup>, with a median of 425 km s<sup>-1</sup>). The lines are especially broad (typically 400 to 600 km s<sup>-1</sup>) in the outflow sources. As shown in Figure 4, the net blueshift in these sources is typically about half the line width ($`\mathrm{\Delta }`$$`v`$ 1/2 $`W`$). The peculiar AGN IRAS11119+3257, with $`W<<\mathrm{\Delta }`$$`v`$, is the notable exception. Thus, in a typical outflow, the redmost absorption occurs close to $`v_{sys}`$. This is strongly suggestive of a flow in which matter is injected at roughly zero velocity and then accelerates outward. The approximate implied terminal velocity of the flow is then $`v_{term}`$ $`\mathrm{\Delta }`$$`v+0.5W`$, which ranges from 220 to 1450 km s<sup>-1</sup> in our sample (Table 4). This picture is quite different from the standard one of a simple expanding ‘superbubble’ in which the absorption is due to a thin layer of cooled post-shock gas, and for which $`W<<\mathrm{\Delta }`$$`v`$ would be expected (e.g. Weaver et al. 1977).
It is instructive to compare the observed velocities in the absorbing material to the velocities expected from purely gravitational forces in the starburst galaxy. This is shown in Figure 5, where we plot $`W`$ vs. the galaxy rotation speed ($`v_{rot}`$ \- see Table 1 for details). This figure has several interesting implications. First, neither the sample as-a-whole nor any of the above subsamples show any correlation between the velocity dispersion in the absorbing material and the galaxy rotation speed. This suggests that gravity does not play a dominant role in determining the dynamics of the absorbing gas.
Figure 5 also shows that the $`NaD`$ lines are surprisingly narrow in the $`SSC`$ sources compared with expectations for either stars or gas in the bulge of the starburst ‘host’ galaxy. The lines are exceptionally narrow if they are stellar in origin, since in this case the observed line broadening ($`W_{obs}`$) will be produced by both the intrinsic stellar line broadening ($`W_{}`$) and that produced by galactic dynamics ($`W_{gal}`$): $`W_{gal}=\sqrt{W_{obs}^2W_{}^2}`$. The observed equivalent widths of the $`Na\lambda `$5890 line are in the range 2.45$`\pm `$0.4 Å in the $`SSC`$ nuclei. If the absorption were purely stellar, the minimum required values for $`W_{}`$ would be 125$`\pm `$20 km s<sup>-1</sup> (corresponding to completely black stellar lines). The typical implied values for $`W_{gal}`$ in the $`SSC`$ sample would then be 60 to 200 km s<sup>-1</sup>, with a median value of 130 km s<sup>-1</sup>.
To emphasize how narrow the lines are in the $`SSC`$ sample, we show in Figure 5 the empirical relation (Whittle 1992; Franx 1993) between the galaxy rotation speed and the bulge velocity dispersion as a function of Hubble type for a sample of normal disk galaxies. The values for $`W_{obs}`$ in the $`SSC`$ objects are on-average $``$ 0.2 dex below this relation for normal galaxies of the same rotation speed and Hubble type (typically Sa to Sc), while the implied values for $`W_{gal}`$ would be even more discrepant (see above). Put another way, based on the Hubble types and the galaxy absolute magnitudes ($`M_B`$ $``$ -19 to -21) for the $`SSC`$ subsample, the Faber- Jackson relation for normal galactic bulges would predict typical values of $`W_{gal}`$ $``$ 200 to 300 km s<sup>-1</sup> (e.g. Nelson & Whittle 1996), while the observed widths are typically only 140 to 200 km s<sup>-1</sup>, even without a correction for the line broadening due to $`W_{}`$.
The nebular emission lines are also narrow in the $`SSC`$ nuclei, as has been shown to be more generally true for starbursts by Weedman (1983). In this case, Lehnert & Heckman (1996b) showed that the narrowness of the nuclear emission-lines could be understood because the ionized gas was rotationally supported and did not fairly sample the galaxy rotation curve (it lies within the region of the galaxy with solid-body rotation). If this explanation applies to the $`NaD`$ lines in the $`SSC`$ nuclei, it implies that a significant fraction of the stellar contribution comes from a dynamically-cold (disk/starburst) component rather than from the bulge.
Finally, Figure 5 shows that the $`NaD`$ linewidths are relatively large in the outflow sources ($`W`$ $``$ 1 to 3 $`v_{rot}`$). As we have argued above, the kinematic properties of the $`NaD`$ profiles suggest that gas is ‘loaded’ into the outflow at $`vv_{sys}`$ and is then accelerated up to some terminal velocity that corresponds to the most-blueshifted part of the $`NaD`$ line profile. We plot $`v_{term}`$ vs. $`v_{rot}`$ in Figure 6, from which it is clear that $`v_{term}`$ is significantly larger than $`v_{rot}`$, but is uncorrelated with it. This suggests that the outflows may be able to selectively escape the shallower galactic potential wells, as we will discuss in section 4.2 below.
Neither the $`SSC`$ nor the $`ISD`$ subsamples show a significant correlation between the widths of the $`NaD`$ absorption-line and the $`H\alpha `$ emission-line. In particular, the outflow sources with very broad (400 to 600 km s<sup>-1</sup>) $`NaD`$ absorption-lines have $`H\alpha `$ emission-line widths ranging from 145 km s<sup>-1</sup> (NGC7552) to 1500 km s<sup>-1</sup> (IRAS11119+3257). This presumably means that the dynamics of the more tenuous outflowing absorbing gas is largely decoupled from that of the dense (high emission-measure) gas that provides most of the nuclear line-emission.
As described above, we have fit the profile of the $`NaD`$ doublet with a single pair of Gaussians constrained to have the same widths and a fixed separation. Inspection of Figure 1 clearly shows that the observed profiles of many of the $`ISD`$ sample are more complex than this. The $`ISD`$ profiles generally have a larger kurtosis than a Gaussian (i.e. narrower core and broader wings) and are sometimes asymmetric with a weak blueward wing on the $`\lambda `$5890 profile (e.g. NGC 1808, IRAS 10565+2448, IRAS 11119+3257, NGC 6240), and/or definite substructure (e.g. NGC 1614, NGC 3256, IRAS 11119+3257). Observations at higher spectral resolution should prove instructive.
### 3.3 The Roles of Luminosity and Geometry
Of the 32 galaxies in our sample, 14 show relatively weak interstellar $`NaD`$ absorption-lines (the SSC sample, in which the stellar contribution to the line is strong), 6 have predominantly interstellar $`NaD`$ lines lying close to the systemic velocity of the galaxy, and 12 have interstellar lines that are blueshifted by more than 100 km s<sup>-1</sup> relative to $`v_{sys}`$. These 12 outflow sources differ systematically from the other objects in two striking respects: they are more luminous starbursts and they are preferentially located in galaxies seen relatively face-on.
Specifically, 64% (9/14) of the galaxies with $`L_{IR}>`$ 10<sup>11</sup> L show outflows, compared to only 28% (5/18) of the less luminous galaxies. The mean values for $`logL_{IR}`$ are 11.44$`\pm `$0.18 and 10.86$`\pm `$0.13 for the outflow and other sources respectively, a difference that is significant at the 2.6 $`\sigma `$ level. The relationship to galaxy inclination is stronger: 69% (11/16) of the galaxies with a ratio of semi-major to semi-minor axes $`a/b`$ 2.0 show outflows, while this is true for only 6% (1/16) of the flatter (more highly inclined) galaxies. The mean values for $`log(a/b)`$ are 0.20$`\pm `$0.03 and 0.42$`\pm `$0.03 for the outflow and other sources respectively, a difference that is significant at the 4.6 $`\sigma `$ level.
It is likely that the primary correlation is between an observed outflow and low galaxy inclination (small $`a/b`$). The weaker apparent correlation with $`L_{IR}`$ is probably induced by the loose anti-correlation in our sample between $`L_{IR}`$ and $`a/b`$. This anti- correlation reflects our selection of galaxies from both the LH95 ‘edge-on’ galaxy sample (large $`a/b`$ and moderate $`L_{IR}`$) and the AHM ‘FIR-warm’ sample (broad range in $`a/b`$ and large $`L_{IR}`$).
Taken at face value, the correlation with galaxy inclination implies that there is a high probability ($``$70%) that an observer located within $``$60 of the rotation axis of a starburst galaxy will see outflowing gas in absorption. This geometrical constraint is consistent with the observed loosely-collimated outflows seen in emission along the minor axes of edge-on starburst galaxies (e.g. Dahlem, Weaver, & Heckman 1998).
### 3.4 Column Densities and Optical Depths
The $`NaD`$ line is clearly optically-thick in these galaxies. The ratio of the equivalent widths of the $`\lambda `$5890 and $`\lambda `$5896 members of the doublet ($`R`$) can be used to estimate the optical depth (e.g. Spitzer 1968). The distribution of $`R`$ is markedly different in the $`SSC`$ and $`ISD`$ subsamples. In the former, there is a very narrow observed range ($`R`$ $``$ 1.1 to 1.3). This is consistent with a strong stellar contribution to the $`NaD`$ line, since $`R`$ $``$ 1.0 to 1.3 (indicative of large optical depths) is characteristic of cool stars. The range is much broader for the $`ISD`$ sample, from $`R`$ = 1.1 to 1.7. This range corresponds to central optical depths in the $`\lambda `$5896 line of $`\tau `$ 20 to 0.5.
At first sight, it might appear odd that the $`NaD`$ line is optically-thick, yet is not black at line center. This can be seen for the $`ISD`$ sample in Figure 7, where we have plotted $`R`$ vs. the normalized residual intensity at the center of the $`\lambda `$5890 feature: $`I_{5890}=F_{5890}/F_{cont}`$ (with the respective fluxes measured at line center and in the adjacent continuum). There is a broad range in $`I_{5890}`$ from 0.14 (nearly black) to 0.7. More tellingly, there is no correlation between $`R`$ and $`I_{5890}`$ . This implies that the absorbing gas does not fully cover the background continuum light, and that $`I_{5890}`$ is determined more by this covering factor ($`C_f`$) than by the optical depth. A covering factor less than unity is natural in these galaxies. First, the continuum light may arise in part from stars in the galaxy that are located in front of most of the absorbing gas (i.e. this is not the idealized case of a purely foreground absorbing screen: the gas and stars are likely to be mixed). Secondly, the gas is likely to be quite clumpy and inhomogeneous (e.g. Calzetti 1997; Gordon, Calzetti, & Witt 1997).
In the limit of large optical depth, $`C_f=(1I_{5890})`$, but for low or moderate optical depth $`C_f>(1I_{5890}`$). For the typical optical depths in this sample, we can approximate $`C_f`$ by (1 - $`I_{5890}`$). This can be demonstrated quantitatively for those members of the $`ISD`$ sample in which the $`NaD`$ lines are well-resolved, narrow enough so that the two doublet members are cleanly separated from one-another ($`W<`$ 300 km s<sup>-1</sup>), and that have high signal-to-noise spectra. These constraints leave us with only three objects: NGC1808, NGC2146, and M82. Following Hamann et al. (1997) and Barlow & Sargent (1997), we have:
$$C_f=(I_{5896}^22I_{5896}+1)/(I_{5890}2I_{5896}+1)$$
(1)
where $`I_{5896}`$ is the normalized intensity at the center of the $`\lambda `$5896 line. The measured values of $`C_f`$ are 0.83, 0.84 and 0.84 for NGC1808, NGC2146, and M82 respectively, while the corresponding values for $`(1I_{5890})`$ are 0.83, 0.82, and 0.82.
In this circumstance - in which optically-thick gas only partially covers the continuum source - the measured equivalent width of the $`NaD`$ doublet ($`EQ`$) will be insensitive to the $`NaI`$ column density, and will instead be primarily determined by the product of $`C_f`$ and the line-of-sight velocity dispersion in the gas. We plot the separate dependences of $`EQ`$ on $`I_{5890}`$ and $`W`$ in Figures 8 and 9 respectively for the $`ISD`$ sample. It is clear from these two figures that $`EQ`$ is determined largely by the covering factor (Figure 8), since there is no correlation between $`W`$ and $`EQ`$ (Figure 9).
Given that the $`NaD`$ doublet is moderately optically-thick in these galaxies, it is not straightforward to estimate a $`NaI`$ column density ($`N_{NaI}`$). We have taken three approaches, and emphasize that these are designed to give us only a rough (order-of-magnitude) estimate. Our techniques can potentially underestimate $`N_{NaI}`$, because they are insensitive to any $`NaI`$ sub-component that is highly optically-thick, yet kinematically quiescent.
The first is the classical doublet ratio method (e.g. Spitzer 1968), which relates $`R`$ directly to the optical depth at line center, and thereby allows the column density to be deduced from the equivalent width. In the spirit of this analysis, we will not attempt to measure columns for all the individual cases, but will instead estimate a characteristic value based on the typical observed parameters. The median value observed in the $`ISD`$ sample is $`R`$ 1.2, implying that the corresponding median optical depth at the center of the $`NaD\lambda `$5896 line is $`\tau _{5896}`$ $``$ 4 (see Table 2.1 in Spitzer 1968). The median observed value $`EQ`$ 6 Å for the doublet, equation 2-41 and Table 2.1 in Spitzer (1968), and the oscillator strength from Morton (1991), together imply $`N_{NaI}`$ $``$ 10<sup>14</sup> cm<sup>-2</sup>. Note that this assumes $`C_f`$ = 1, and should be increased by $`C_f^1`$, or a typical factor of $``$ 1.6.
A variant of the doublet-ratio technique can be applied to the three cases discussed above in which the two members of the $`NaD`$ doublet are cleanly separated and well-resolved (NGC1808, NGC2146, and M82). Again, following Hamann et al (1997) we have:
$$\tau _{5896}=ln[C_f/(I_{5896}+C_f1)]$$
(2)
The resulting values for $`\tau _{5896}`$ are 2.3, 2.1, and 1.9 for NGC1808, NGC2146, and M82 respectively. These are smaller than the values implied by $`R`$ by a factor of $``$ 2 in these cases. The implied values for $`N_{NaI}`$ are 1.0 $`\times `$ 10<sup>14</sup>, 6 $`\times `$ 10<sup>13</sup>, and 6 $`\times 10^{13}`$ cm<sup>-2</sup> after correction by $`C_f^1`$.
We have also measured the equivalent width of the $`KI\lambda `$7665 line in three of the nuclei (NGC1614, NGC1808, and NGC3256). Since $`KI`$ and $`NaI`$ have very similar ionization potentials, and since $`K`$ and $`Na`$ show similar grain depletion patterns (Savage & Sembach 1996), the expected ratio of the $`NaI`$ and $`KI`$ column densities should be 15 for gas with a solar Na/K ratio. The measured values for the $`NaD`$ doublet ratio imply optical depths at the center of the $`NaD\lambda `$5890 line of 8, 16, and 1.6 for NGC1614, NGC1808, and NGC3256 respectively. The implied optical depths for the $`KI\lambda `$7665 line would then be 0.5, 1.1, and 0.1 respectively. Using the oscillator strength tabulated by Morton (1991), the measured equivalent widths of the line imply that $`N_{KI}`$ = 3 $`\times `$ 10<sup>12</sup>, 4 $`\times `$ 10<sup>12</sup>, and 1.3 $`\times `$ 10<sup>12</sup> cm<sup>-2</sup> respectively. Assuming $`N_{NaI}`$ = 15 $`N_{KI}`$, the corresponding $`NaI`$ columns are 4.5 $`\times `$ 10<sup>13</sup>, 6 $`\times `$ 10<sup>13</sup>, and 2 $`\times `$ 10<sup>13</sup> cm<sup>-2</sup>. These values are about a factor of two or three smaller than would have been deduced for these three cases using the $`NaD`$ doublet ratio alone. The value for NGC1808 is in good agreement with that derived from Equation 2. Under the circumstances, we regard the agreement between the three methods as satisfactory, and conclude that the typical value in the $`ISD`$ sample is $`logN_{NaI}`$ = 13.5 to 14. A final indirect indication that these $`NaI`$ column densities are roughly correct comes from the detections of the “Diffuse Interstellar Bands” in the seven highest-quality spectra of the $`ISD`$ sample (Heckman & Lehnert 2000). The observed strengths of these features agree with the strengths seen in Galactic sight-lines with $`logN_{NaI}`$ 13.5 to 14.
What is the total gas column density associated with the outflow? To calculate this directly from the (already uncertain) $`NaI`$ column requires knowing the metallicity of the gas, the fractional depletion of $`Na`$ onto grains (typically a factor of $`10`$ in diffuse clouds in the Milky Way) and the potentially substantial ionization correction to account for ionized $`Na`$. Assuming solar $`Na`$ abundances and a factor of ten correction for depletion onto grains (e.g Savage & Sembach 1996), $`N_{NaI}`$ = 10<sup>14</sup> cm<sup>-2</sup> implies a typical value for $`N_H`$ of 5 $`\times `$ 10<sup>20</sup> ($`N_{Na}/N_{NaI}`$) cm<sup>-2</sup>. We can also take an empirical approach suggested by the correlation between $`N_{NaI}`$ and the total gas column towards stars in our own Galaxy. Using the data in Herbig (1993), values for $`N_{NaI}`$ in the range we estimate ($`logN_{NaI}`$ 13.5 to 14) correspond to sight-lines with $`N_H`$ $``$ 1.5 to 4 $`\times `$ 10<sup>21</sup> cm<sup>-2</sup>. Interestingly, this is just the range of values for $`N_H`$ deduced from the amount of reddening along the line-of-sight to these nuclei based on either the Balmer decrement or the colors of the optical continuum, assuming a normal Galactic extinction-curve and dust-to gas ratio ( section 3.5, and see also AHM; Veilleux et al 1995).
These estimates suggest that the ionization correction factor is significant but not huge (i.e. $`N_{Na}/N_{NaI}`$ 3 to 10). Since its ionization potential is only 5.1 eV, the presence of relatively significant amounts of $`NaI`$ implies that it is associated with gas having a significant dust optical depth in the near-UV: for a Galactic extinction curve and dust-to-gas ratio, a Hydrogen column density of $`N_H=8\times 10^{20}`$ cm<sup>-2</sup> is required to produce $`\tau _{dust}=1`$ at 5.1 eV ($`\lambda `$ 2420 Å).
It is instructive to compare the total column densities we infer for the outflows of a few $`\times `$ 10<sup>21</sup> cm<sup>-2</sup> to the column densities in the other components of the ISM in these galaxies. Column densities to the nucleus for the hot X-ray-emitting gas are estimated to be of-order 10<sup>21</sup> cm<sup>-2</sup> in superwinds (e.g. Suchkov et al 1994; Heckman et al 1999; Strickland 1998). In the nuclei themselves, the dominant ISM component is molecular, and the inferred columns range from $``$ 10<sup>23</sup> to 10<sup>25</sup> cm<sup>-2</sup> (e.g. Sanders & Mirabel 1996).
The $`HI\lambda `$21cm line is observed in absorption against the bright nonthermal radio sources in starburst nuclei (e.g. Koribalski 1996; Heckman et al 1983; Mirabel & Sanders 1988). The implied column densities are typically a few $`\times `$ 10<sup>21</sup> to 10$`{}_{}{}^{22}(T_{spin}/100K)`$ cm<sup>-2</sup>. The absorption is centered close to $`v_{sys}`$ and spans a velocity range similar to that of the molecular gas. This strongly suggests that this gas is a trace atomic component in the starbursting molecular disk or ring. The kinematics of the gas responsible for the $`\lambda `$21cm absorption are therefore quite distinct from the gas that produces the blueshifted $`NaD`$ absorption. This has several plausible explanations. First, the outflowing $`HI`$ is probably too hot ($`T>10^3`$ K) to produce strong absorption at $`\lambda `$21cm. Second, the background radio continuum source against which the gas that produces the $`\lambda `$ 21cm absorption is observed will almost certainly be invisible in the optical: it lies behind a total column density (overwhelmingly $`H_2`$) of $``$ 10<sup>23</sup> to 10<sup>25</sup> cm<sup>-2</sup>, corresponding to $`A_V`$ = 60 to 6000! Clearly, such material will not contribute to the observed $`NaD`$ absorption-lines.
### 3.5 Dust Associated with the Absorbing Gas
We have argued in section 3.4 that the $`NaD`$ lines are are optically thick, and that $`EQ`$ is set primarily by the covering fraction for the absorbing gas ($`C_f`$) rather than by the line width ($`W`$ \- see Figures 8 and 9). This inference helps explain the otherwise puzzling correlations found by AHM and Veilleux et al (1995) between $`EQ`$ and the reddening inferred from either the Balmer decrement or the color of the optical continuum in the nuclear spectra of large samples of starbursts. For $`\tau _{NaD}`$ $`>>`$ 1 and $`C_f`$ = 1, $`EQ`$ would be set by $`W`$, and so no correlation with the reddening would be expected. If instead $`EQ`$ is principally determined by the fraction of the starburst that is covered by gas containing $`NaI`$ (and dust grains), then this correlation is more reasonable.
Veilleux et al (1995) have also shown (via the Balmer decrement) that the region of significant reddening extends far beyond the nucleus in many far-IR-bright galaxies. Thus, to gain further insight into the relationship between the $`NaD`$ absorption and dust-reddening, we have mapped out the spatial variation in the depth of the $`NaD`$ line ($`I_{5890}`$) and the reddening in the six galaxies in our $`ISD`$ sub-sample for which we have the relevant data on the reddening (M82, NGC3256, NGC6240, Mrk273, IRAS03514+1546, and IRAS10565+2448). The size of the region mapped was set by the detectability of the $`NaD`$ line, and ranges from 3 to 9 kpc (except for M 82, where the mapped region is only 500 pc in diameter). In each case, we have corrected the H$`\alpha `$ and H$`\beta `$ emission-line fluxes for the effects of stellar absorption-lines (using measures of the equivalent widths of the high-order stellar Balmer absorption-lines in NGC 3256 and NGC 6240 and an assumed value of 2 Å for the other galaxies). We have also corrected the data for foreground reddening using the measured Galactic $`HI`$ column density and assuming a standard extinction curve.
Figure 10 shows that not only do the extra-nuclear data points for these six galaxies define a good correlation between the amount of reddening and the depth of the $`NaD`$ line along a given line-of-sight through the starburst and its outflow, they define the same correlation as that defined by the ensemble of all the $`ISD`$ nuclei in our sample. The nuclear and off-nuclear points are pretty well-mixed in Figure 10, although there is some tendency for the nuclear lines-of-sight to have the larger values of reddening and deeper $`NaD`$ absorption-lines. The correlation of $`I_{5890}`$ is better with the color of the stellar continuum than with the Balmer decrement. This is reasonable because the $`NaD`$ line is observed in absorption against the background stellar continuum (rather than against the emission-line gas) and because the Balmer decrement is likely to be significantly affected by dust directly associated with the emission-line gas itself (in addition to the dust in the foreground material responsible for the $`NaD`$ absorption).
The observed Balmer decrements imply extinctions of $`A_V`$ $``$ 1 to 5 for a standard Galactic extinction curve. Similar values are implied by the continuum colors: a typical starburst is predicted (in the absence of reddening) to have a color of $`log[C_{65}/C_{48}]`$ -0.3 (Leitherer & Heckman 1995), while the observed colors in Figure 10 range from $`log[C_{65}/C_{48}]`$ -0.2 to +0.3 (corresponding to $`A_V`$ = 0.7 to 4.2). Note also that in both Figure 10a and 10b, the extrapolation of the correlation to $`I_{5890}`$ = 1.0 (no absorption, $`C_f`$ = 0) has an x-intercept at the intrinsic values expected for an unreddened starburst ($`log[C_{65}/C_{48}]`$ -0.3 and $`log[H\alpha /H\beta ]`$ = 0.46).
In summary, the data imply that over regions with sizes of several or many kpc, the outflows contain inhomogeneous highly dusty material. For a standard Galactic extinction law and dust-to-gas ratio, the typical implied $`HI`$ columns are a few $`\times `$ 10<sup>21</sup> cm<sup>-2</sup>. These $`HI`$ column densities agree well with the estimates in section 3.4 above based upon the $`NaI`$ column density.
### 3.6 Sizes, Masses, and Energies
We have measured the size of the region over which significantly blueshifted $`NaD`$ absorption is detected ($`\mathrm{\Delta }`$$`v`$ $`>`$ 100 km s<sup>-1</sup>) for the 12 outflow sources. The sizes are listed in Table 4, and range from 1 to 10 kpc in diameter. They must be regarded as lower limits (since the background starlight usually becomes too faint to detect the absorption at larger radii). Tracing the full extent of the absorbing material farther out into the galactic halos will probably require observing suitably bright background QSO’s (see Norman et al 1996).
These lower limits to the size of absorbing region can be used to estimate the (minimum) mass and kinetic energy in the outflow. That is, for a region with a surface area $`A`$, a column density $`N_H`$, and an outflow velocity $`\mathrm{\Delta }`$$`v`$:
$$M>5\times 10^8(A/10kpc^2)(N_H/3\times 10^{21}cm^2)M_{}$$
(3)
$$E>2\times 10^{56}(A/10kpc^2)(N_H/3\times 10^{21}cm^2)(\mathrm{\Delta }v/200km/s)erg$$
(4)
We have scaled these relations using values for $`A`$, $`N_H`$, and $`\mathrm{\Delta }`$$`v`$ that are typical, and have assumed an equal contribution to $`M`$ and $`E`$ from the front (observed) and back sides of the outflow.
If we adopt a simple model of a constant-velocity, mass-conserving superwind flowing into a solid angle $`\mathrm{\Omega }_w`$, extending to arbitrarily large radii from some minimum radius ($`r_{}`$ \- taken to be the radius of the starburst within which the flow originates), we obtain:
$$\dot{M}60(r_{}/kpc)(N_H/3\times 10^{21}cm^2)(\mathrm{\Delta }v/200km/s)(\mathrm{\Omega }_w/4\pi )M_{}/yr$$
(5)
$$\dot{E}8\times 10^{41}(r_{}/kpc)(N_H/3\times 10^{21}cm^2)(\mathrm{\Delta }v/200km/s)^3(\mathrm{\Omega }_w/4\pi )erg/s$$
(6)
The statistics of the $`ISD`$ subsample in the present paper imply that outflows are commonly observed in absorption in IR-selected starbursts (12/18 cases). On the other hand, many of the outflow galaxies in the present sample were selected from AHM on the basis of the strength of their $`NaD`$ line (objects above the $``$ 70th percentile in $`EQ`$). If the presence of observable blueshifted absorption is determined by viewing angle (see section 3.3), this suggests that $`\mathrm{\Omega }_w`$/4$`\pi `$ lies in the range $``$ 0.2 to 0.6 (consistent with the weakly-collimated bipolar outflows seen in well-studied superwinds).
To put the above estimates into context, we can consider the rate at which mass and energy are returned by massive stars. The median bolometric luminosity of the 12 outflow galaxies in our sample is $`L_{bol}`$ 2 $`\times `$ 10<sup>11</sup> L. The implied median rates of mass and kinetic energy returned from supernovae and stellar winds are roughly $`\dot{M_{ret}}`$ = 5 M per year and $`\dot{E_{ret}}`$ = 10<sup>43</sup> erg s<sup>-1</sup> respectively (e.g. Leitherer & Heckman 1995). Since $`(\dot{M}/\dot{M_{ret}})`$ 3 to 10, the absorption-line gas in the outflow must be primarily ambient gas that has been loaded into the flow. This inference agrees with similar conclusions about the hot X-ray-emitting gas in superwinds (section 4.2 below, and see e.g. Strickland 1998; Suchkov et al 1996; Heckman et al 1999). Since $`(\dot{E}/\dot{E_{ret}})<`$ 10%, the absorbing gas does not carry the bulk of the energy supplied by the starburst. Most of this energy probably resides in the form of the thermal and kinetic energy of the much hotter ($`T>10^{5.5}`$ K) X-ray-emitting gas.
A bolometric luminosity of 2 $`\times `$ 10<sup>11</sup> L corresponds to a star-formation rate of about 12 M per year (for a Salpeter IMF extending from 1 to 100 M). Thus, the outflow rates estimated from the $`NaD`$ lines are comparable to the star-formation rate: the feedback from massive stars drives the ejection of as much gas as is being converted into stars. Similar inferences for starbursts have been made using the X-ray and optical emission-line data (e.g. Suchkov et al 1996; Heckman et al 1999; Della Ceca et al 1996,1999; Martin 1999).
## 4 Discussion & Implications
### 4.1 The Origin & Dynamics of the Absorbing Material
As discussed above, the red-most part of absorption-line profile in the outflow objects is close to $`v_{sys}`$, suggesting that absorbing material is injected from quiescent material at or near $`v_{sys}`$, and is then accelerated up to some terminal velocity as it flows outward. This is physically plausible, as the hot (X-ray emitting) outflowing gas interacts hydrodynamically with colder denser material that is located either inside the starburst, or in the inner portions of the galactic halo (see for example Hartquist, Dyson, & Williams 1997; Suchkov et al 1994; Strickland 1998).
Let us assume that a cloud of gas with a column density $`N`$, originally located a distance $`r_0`$ from the starburst, is accelerated by a constant- velocity superwind that carries an outward momentum flux $`\dot{p}`$ into a solid angle $`\mathrm{\Omega }_w`$. Ignoring the effects of gravity for the moment, the clump’s terminal velocity will be (Strel’nitskii & Sunyaev 1973):
$$v_{term}=420(\dot{p}/7\times 10^{34}dynes)^{1/2}(\mathrm{\Omega }_w/1.6\pi )^{1/2}(r_0/kpc)^{1/2}(N/3\times 10^{21}cm^2)^{1/2}km/s$$
(7)
In this expression, we have used the momentum flux supplied by stellar winds and supernovae (Leitherer & Heckman 1995) in a starburst having a bolometric luminosity equal to the median value for the outflow sample (2 $`\times 10^{11}`$ L) and have adopted the estimates in section 3 above for $`N`$ and $`\mathrm{\Omega }_w`$. Starbursts with this luminosity have typical estimated radii of roughly 1 kpc (see for example Heckman, Armus, & Miley 1990; Meurer et al 1997). From these elementary considerations, we conclude that the observed terminal velocities (typically 400 to 600 km s<sup>-1</sup>) are easily accommodated.
Equation 7 also predicts that the outflow speeds will be larger in more luminous starbursts. This trend is mitigated to some degree by the fact that more powerful starbursts tend to be larger. Lehnert & Heckman (1996b) and Meurer et al (1997) argue that starbursts have a maximum characteristic surface-brightness, which then implies $`\dot{p}L_{bol}r^2`$. Together with Equation 7, this implies that such ‘maximum starbursts’ will have $`v_{term}L_{bol}^{1/4}`$ (although the clouds can not be accelerated to velocities larger than that of the flow that accelerates them!). Our sample shows no convincing evidence of a trend for larger $`v_{term}`$ in the more luminous systems, but this sample covers a rather small range in starburst luminosity. It will be instructive to extend this study to dwarf starbursts with $`L_{bol}<10^9`$ $`L_{}`$.
Assume now that the cloud immersed in the outflow is subjected to a gravitational force imposed by an isothermal galaxy potential whose depth corresponds to a circular rotation speed $`v_{rot}`$. In order that the outwardly-directed force due to the superwind exceed the inwardly-directed force of gravity, the value for the cloud column density must satisfy the condition:
$$N<7\times 10^{21}(\dot{p}/7\times 10^{34})(\mathrm{\Omega }_w/1.6\pi )^1(r/kpc)^1(v_{rot}/200km/s)^2cm^2$$
(8)
Thus, the typical column densities estimated for the outflows ($``$ few $`\times `$ 10<sup>21</sup> cm<sup>-2</sup>) lie near the upper bound for material that will flow out (rather than falling in). This may not be a coincidence: given a range of cloud column densities, the blueshifted absorption-line will be dominated by the largest-column-density clouds that can be expelled. Alternatively, the observed column densities may simply arise because $`N_H2\times `$ 10<sup>21</sup> cm<sup>-2</sup> corresponds to a dust optical depth of unity in the continuum at the wavelength of the $`NaD`$ doublet (for a standard Galactic dust- gas ratio). That is, continuum-emitting regions in these nuclei lying behind sight-lines with much higher columns are invisible in optical light and sources behind sight-lines with much lower columns contain little $`NaI`$. It would be interesting to test Equation 8 by measuring values of $`N`$ for outflows in ‘low-intensity’ starbursts (small $`\dot{p}/r`$) and starbursts occurring in dwarf galaxies (small $`v_{rot}`$).
### 4.2 Insights from Numerical Simulations
The above elementary considerations give some simple physical insights into the origin and dynamics of the absorbing material. More detailed insight comes from hydrodynamical simulations of starburst-driven superwinds (e.g. Tomisaka & Bregman 1993, Suchkov et al 1994, Strickland 1998, Tenorio-Tagle & Munoz-Tunon 1998). In these simulations the coolest densest gas that has been hydrodynamically disturbed by the starburst is associated with the swept-up shell of ISM that propagates laterally in the plane of the galaxy, and fragments of the cap of the original superbubble shell now being carried vertically out of the disk by faster, more tenuous, wind material.
Shear between the hot shocked starburst ejecta and the cool dense shell in the disk of the galaxy leads to entrainment and stripping of cool dense gas into the wind flowing out of the disk (through Kelvin-Helmholtz instabilities, and presumably additional interchange processes such as thermal conduction and turbulent mixing layers that can not be included in current simulations). Dense gas already in the wind interior, for example the superbubble shell fragments, is accelerated outward via the ram pressure of the wind.
This process can be seen in Figure 11, in which the outward trajectories of four typical entrained clouds are traced over an interval of 1.5 Myr from a 2-D hydrodynamic simulation of M82’s galactic wind (Strickland 1998). This is based on the thick-disk ISM distribution of Tomisaka & Bregman (1993). In this model a mass of 10<sup>8</sup> M is turned into stars in an instantaneous burst (Salpeter IMF over the mass range 1 to 100 M). At a time 7 Myr after the burst, the resulting wind has properties that are a reasonable match to M82. By this time, supernovae and stellar winds have returned 1.3 $`\times `$ 10<sup>7</sup> M and 6 $`\times `$ 10<sup>56</sup> ergs to the ISM.
Within $`|z|1.5`$ kpc of the disk there is $`M=1.9\times 10^8M_{}`$ of gas cooler than $`T=3\times 10^5`$ K. The majority of this gas is at the minimum temperature allowed in these simulations of $`T6\times 10^4`$ K. This material occupies a projected area of $`2`$ kpc<sup>2</sup>, so the average hydrogen column density of this gas is $`N_\mathrm{H}9\times 10^{21}`$ cm<sup>-2</sup>. This cool gas has a broad range of velocities, from $`v10`$$`10^3`$ km s<sup>-1</sup>, with a mode of $`60`$ km s<sup>-1</sup> (which is associated with the slow expansion of the outer shock in the plane of the galaxy). The associated kinetic energy is 1.3 $`\times `$ 10<sup>55</sup> ergs.
At higher distances above the disk there is much less cool gas. For gas above $`|z|=1.5`$ kpc the mass of cool gas, associated primarily with the fragments of superbubble shell cap, is $`M=1.5\times 10^7M_{}`$. For a projected area of $``$ 8.5 kpc<sup>2</sup>, the average column density is $`N_\mathrm{H}1.6\times 10^{20}`$ cm<sup>-2</sup>. This cool gas high above the plane typically has higher velocities than the gas within the plane of the galaxy, the distribution of mass with velocity being approximately flat between velocities of $`v=10^2`$$`10^3`$ km s<sup>-1</sup>. The kinetic energy associated with this gas is 1.7 $`\times `$ 10<sup>55</sup> ergs.
Thus, the total mass of cool gas in the wind is 2 $`\times `$ 10<sup>8</sup> M. This is twice as large as the mass of stars formed in the burst, and is 16 times larger than the mass directly returned by massive stars. In contrast, the total kinetic energy in the cool gas (3 $`\times `$ 10<sup>55</sup>) is only 5% of the kinetic energy returned by massive stars. The lion’s share of this energy is the form of thermal and kinetic energy of the hot ($`T>10^{5.5}`$ K) gas in the wind.
It is worth noting several of the limitations of these simulations with respect to their treatment of the cool dense ISM: none of these simulations can explicitly include the cool dense clouds of material within the ISM and starburst region that are known to exist, and are thought to play a key role in “mass-loading” the outflow (e.g. Hartquist, Dyson, & Williams 1997). As a result, the cool dense gas in these simulations is confined to larger radii near the outer shock and to the shell fragments (e.g. there is no way of explicitly treating the entrainment of clouds from within the starburst region itself). <sup>2</sup><sup>2</sup>2The entrainment of material into the hotter, more tenuous phases, from the hydrodynamical destruction of such clouds can and has been simulated, but these “mass-loaded” simulations (Suchkov et al 1996, Strickland 1998) do not consistently treat the properties of the clouds themselves. These simulations also have a minimum allowed gas temperature of $`T6\times 10^4\mathrm{K}`$, due to the method of simulating ISM turbulent pressure support by an enhanced thermal pressure. Hence all the gas that would in reality have lower temperature is forced to have this minimum temperature, which in turn affects the density of this gas, and prevents us from knowing the exact distribution of this mass between the gas phases cooler than this minimum temperature. Similarly the processes of entrainment and acceleration of cool gas into the wind are both uncertain and occur at (or below) the scale of the physical resolution of these simulations.
Nevertheless, the results of these simulations are encouragingly similar to the observed parameters in our sample of starburst outflows: the cool gas is predicted to have column densities of several $`\times 10^{20}`$ to $`10^{22}\mathrm{cm}^2`$, a mass that is comparable to that of the stars formed in the burst, outflow velocities in the range $`v10^2`$$`10^3\mathrm{km}\mathrm{s}^1`$, and a kinetic energy that is of-order 10<sup>-1</sup> of the total kinetic energy returned to the ISM by the starburst. Since the cool gas was originally cold dense material entrained and accelerated by the hot outflow, the presence of substantial amounts of dust (section 3.5) is perhaps not too surprising.
### 4.3 The Fate of the Outflow & the Chemical Evolution of Galaxies
As Figure 6 shows, the inferred terminal velocity in the outflows is typically 400 to 600 km s<sup>-1</sup>, or about two to three times larger than the rotation speed of the starburst’s host galaxy. Are these velocities sufficient to expel the gas from the galaxy altogether, or will the gas return to the galactic disk as a fountain flow?
For an isothermal gravitational potential that extends to a maximum radius $`r_{max}`$, and has a virial velocity $`v_{rot}`$, the escape velocity at a radius $`r`$ is given by:
$$v_{esc}=[2v_{rot}(1+ln(r_{max}/r))]^{1/2}$$
(9)
Thus, $`v_{esc}`$ = 3.0 $`v_{rot}`$ for $`(r_{max}/r)`$ = 33 (e.g. $`r`$ = 3 kpc and $`r_{max}`$ = 100 kpc). As shown in Figure 6, the estimated terminal velocities in the outflows are typically $`v_{term}`$ 2 $`v_{rot}`$, but $`v_{term}`$ is uncorrelated with $`v_{rot}`$.
Similar results have been obtained for the hot X-ray-emitting gas in starburst galaxies. This gas has temperatures of a few to ten million K in dwarf galaxies (e.g. Della Ceca et al 1996; Strickland, Ponman, & Stevens 1997), $`L_{}`$ disk galaxies (e.g. Dahlem, Weaver, & Heckman 1998; Read, Ponman, & Strickland 1997), and extremely powerful starbursts in galactic mergers (e.g. Heckman et al 1999; Moran, Lehnert, & Helfand 1999; Read & Ponman 1998). Martin (1999) has used these X-ray data to estimate that the gas will escape from galaxies with $`v_{rot}<`$ 130 km s<sup>-1</sup>.
We can place these disparate data on common ground by comparing the kinetic ($`NaD`$) and thermal (X-ray) energy per particle to the energy needed for escape. For convenience, we do so by defining an energetically- equivalent velocity for the X-ray gas. The terminal velocity in an adiabatic superwind fed by gas at a temperature $`T_X`$ will be $`v_X(5kT_X/\mu )^{1/2}`$, where $`\mu `$ is the mean mass per particle (Chevalier & Clegg 1985). <sup>3</sup><sup>3</sup>3 This is a conservative approach as it ignores any kinetic energy the X-ray-emitting gas may already have. Currently the velocity and kinetic energy of the X-ray-emitting material in superwinds can not be measured directly, but numerical simulations suggest that the kinetic energy of the hot gas is typically 2 to 3 times its thermal energy (Strickland 1998).
These results are shown in Figure 12, which includes the data from Fig. 6 plus 14 far-IR-bright galaxies for which analyses of broad-band ($``$ 0.1 to 10 keV) X-ray data have been published. These are: M82, NGC253, NGC3628, NGC3079, NGC4631 (Dahlem, Weaver, & Heckman 1998), NGC1569 (Della Ceca et al. 1996), NGC1808 (Awaki et al. 1996), NGC2146 (Della Ceca et al. 1999), NGC3256 (Moran, Helfand, & Lehnert 1999), NGC3310 (Zezas, Georgantopoulos, & Ward 1998), NGC4038/4039 ( Sansom et al. 1996), NGC4449 (Della Ceca, Griffiths, & Heckman 1997), NGC6240 (Iwasawa & Comastri 1998), and Arp299 (Heckman et al 1999). In the cases where two-temperature plasma models were fit to the X-ray data, we have plotted both the corresponding outflow velocities. The agreement between the two data sets is satisfactory. There are three members of the $`NaD`$ outflow sample with X-ray data in Figure 12, and the agreement between the $`NaD`$ terminal velocity and the X- ray temperatures is reasonably good: $`v_{term}`$ = 700 km s<sup>-1</sup> vs. $`v_X`$ = 520 and 780 km s<sup>-1</sup> for NGC1808, $`v_{term}`$ = 580 km s<sup>-1</sup> vs. $`v_X`$ = 490 and 800 km s<sup>-1</sup> for NGC3256, and $`v_{term}`$ = 580 km s<sup>-1</sup> vs. $`v_X`$ = 700 and 940 km s<sup>-1</sup> for NGC6240. This suggests that the fastest-moving $`NaD`$ absorbers are roughly co-moving with the hot superwind fluid.
Figure 12 strongly suggests that shallower galaxy potential wells will be less able to retain the newly-synthesized metals that are returned to the ISM in the aftermath of a starburst. As has been suggested many times (e.g. Wyse & Silk 1985; Lynden-Bell 1992; Kauffmann & Charlot 1998) the selective loss of metal-enriched gas from shallower potential wells could explain both the mass-metallicity relation and radial metallicity gradients in elliptical galaxies and galaxy bulges (Bender, Burstein,& Faber 1993; Franx & Illingworth 1990; Carollo & Danziger 1994; Jablonka, Martin, & Arimoto 1996; Pahre, de Carvallo, & Djorgovski 1998; Trager et al 1998).
A simple prediction of this idea would be that the relationship between metallicity and escape velocity should saturate (flatten) for the deepest potential wells - i.e. locations where the local escape velocity exceeds the velocity of the outflowing metal-enriched gas. Lynden-Bell (1992) has parameterized this in a simple physically-motivated fashion by positing that the fraction of metals produced by massive stars that are retained by the galaxy ($`f_{retained}`$) is proportional to the depth of the galaxy‘s potential well ($`\mathrm{\Phi }`$) for low-mass galaxies, but asymtotes to $`f_{retained}`$ = 1 for the most massive galaxies. We chose to cast his formulation as follows:
$$f_{retained}=v_{esc}^2/(v_{esc}^2+v_{term}^2)$$
(10)
Here $`v_{term}`$ is some characteristic velocity associated with the mixture of supernova (and stellar wind) debris and entrained gas that is ejected from the starburst. It is assumed to be a constant. For low-mass galaxies with $`v_{esc}<<v_{term}`$, $`f_{retained}v_{esc}^2\mathrm{\Phi }`$, or $`f_{retained}L_{gal}^{1/2}`$ via the Faber-Jackson relation. Lynden-Bell shows that this simple formula can reproduce the observed mass-metallicity relation for elliptical galaxies over a range of $``$ 10<sup>6</sup> in galaxy mass, and finds that the characteristic mass at which $`f_{retained}=1/2`$ (e.g. a galaxy in which $`v_{esc}=v_{term}`$) corresponds to an elliptical with $`M_B`$ $``$ -18 (adjusted to our assumed value of $`H_0`$ = 70). Such a galaxy would have a line-of-sight velocity dispersion $`\sigma `$ 140 km s<sup>-1</sup>, corresponding to $`v_{rot}=\sqrt{2}\sigma `$ 200 km s<sup>-1</sup> (Binney & Tremaine 1987). Using equation 9 above, this would correspond to $`v_{esc}`$ 600 km s<sup>-1</sup>. This in turn is a reassuringly good match to the characteristic superwind outflow speeds implied by Figure 12 ($``$ 400 to 800 km s<sup>-1</sup>).
Thus, starburst-driven outflows might imprint a relationship between metallicity and mass in ellipticals (and bulges) over most of observed ranges for these two parameters. While the loss of metal-enriched gas has the most severe impact on dwarf elliptical galaxies, it may nevertheless have general significance in galaxy chemical evolution.
### 4.4 The Metal-Enrichment of the Intergalactic Medium
The data discussed in this paper directly establish the flow of metals out of highly-actively-star-forming galaxies in the local universe, and the process is observed at high-redshift as well (Franx et al 1997; Pettini et al 1998, 1999). Such data allow us to estimate the column densities, outflow rates, and outflow speeds of this material as a function of the rate of star-formation. Meanwhile, over the past few years, the rate of high-mass star-formation over the history of the universe has been measured for the first time (e.g. Madau et al 1996; Steidel et al 1999; Barger et al 1999). This emboldens us to attempt to estimate the amount of metals that have flowed out of galaxies, thereby polluting the inter-galactic medium, over the course of cosmic time.
The discussion in section 3.4 above implies that gas is flowing out of starbursts at a rate that is proportional to the rate of star-formation: $`\dot{M}=\alpha \dot{M_{}}`$ where $`\alpha `$ is one-to-a-few (see also Martin 1999). The present-day mass in stars will be smaller than the total mass turned into stars, since mass has been subsequently returned from these stars: $`M_{,0}=\beta M_{}`$ where $`\beta `$ 0.7 is reasonable for an old present-day system like an elliptical or bulge. The discussion in section 4.3 implies that the outflowing gas will be mostly retained by galaxies with the deepest potential wells, but mostly lost by the less massive systems. Integrating equation 10 above over a Schechter luminosity function implies that $`\dot{M_{lost}}=\gamma \dot{M}`$ with $`\gamma `$ 0.5 (depending on the value of $`v_{term}`$ and the ‘mapping’ of $`M_B`$ to $`v_{esc}`$ in spheroids).
We then assume that over cosmic time, we can attribute the construction of spheroidal systems (elliptical and bulges) to starbursts (see Kormendy & Sanders 1992; Elmegreen 1999; Renzini 1999; Lilly et al. 1999). If we further assume that all star-formation in spheroids over the history of the universe ejected gas at the relative rate seen in local starbursts, then the ratio of the mass of lost-gas to present-day stars in spheroids would be $`\alpha \gamma /\beta `$ (of-order unity). Fukugita, Hogan, & Peebles (1998 - hereafter FHP) estimate that the stars in spheroidal systems today comprise $`\mathrm{\Omega }_{,sph}`$ = 2.6 $`\times `$ 10<sup>-3</sup> (for $`H_0`$ = 70). Since the implied value for the gas expelled from forming spheroids is comparable to this, this gas is therefore a significant repository of baryons, but only of-order 10<sup>-1</sup> of the total estimated baryonic content of the universe (FHP). In rich clusters, nearly the entire stellar mass resides in spheroidal systems, while the cluster potential well is deep enough to have retained all the mass expelled by superwinds (e.g. Renzini 1997). The observed average ratio of the mass of the intra-cluster medium to the stellar mass is $``$ 6 (FHP), so gas ejected by superwinds during spheroid formation would comprise a significant, but minority share of this.
Equation 10 also implies that - integrated over the spheroid luminosity function - roughly half of the metals produced by the stars will have been lost from the galaxies and reside in the intergalactic medium or intracluster medium. Once mixed with the metal-poor “primordial” baryons, the net metallicity would be $``$ 1/6th solar in both the intracluster medium and the general inter-galactic medium (assuming the FHP global value $`\mathrm{\Omega }_{,sph}/\mathrm{\Omega }_{IGM}`$ 6, and assuming a mass-weighted mean metallicity equal to solar for stars in spheroids). The estimated metallicity agrees reasonably well with the measured value of 0.3 solar in rich clusters (e.g. Renzini 1997). A measure of the metal content of the present-day general IGM may be possible with the next generation of UV and X-ray space spectrographs (e.g. Cen & Ostriker 1999).
These are not new arguments by any means (e.g. Gibson, Loewenstein, & Mushotzky 1997; Renzini 1997). What is new is that we are now in a position to observationally verify that intense starbursts of the kind that plausibly built galactic spheroids do indeed drive mass and metals out at a rate and velocity perhaps high enough to account for the observed inter-galactic metals. The presence of such substantial amounts of inter-galactic metals does not violate constraints imposed by the “Madau diagram” (star-formation rate vs. redshift), once reasonable corrections for the effects of dust-extinction are made, nor does it violate the limits set by the far-IR/sub-mm cosmic background (see for example Calzetti & Heckman 1999; Renzini 1997).
### 4.5 The Outflow of Dust
The strong correlation between reddening and the strength of the $`NaD`$ line in starbursts (AHM; Veilleux et al 1995; Figure 10) implies that there is an intimate relationship between the dust and gas, especially given the close way in which the two track one another spatially throughout the outflow (section 3.5, Figure 10, and see Phillips 1993 for the spectacular case of NGC 1808). Moreover, we have argued above that significant dust column densities in the absorbing matter are needed to shield the $`NaI`$ from photoionization by the starburst’s intense UV radiation field.
We therfore conclude that dust is being expelled from starbursts at a significant rate. More quantitatively, for normal Galactic dust, the observed reddening implies a dust surface mass density in the outflow region of $``$10<sup>-4</sup> gm cm<sup>-2</sup>, an outflowing dust mass of $``$ 10<sup>6</sup> to 10<sup>7</sup> M (see equation 3), and a dust outflow rate of 0.1 to 1 M yr<sup>-1</sup> (see equation 5).
Additional evidence for dusty galactic outflows comes from a variety of observations. Spectroscopy with $`HST`$ and $`HUT`$ has established that - just as in the case of the $`NaD`$ lines - the strong $`UV`$ interstellar absorption-lines are frequently blueshifted by several hundred km s<sup>-1</sup> in local starbursts (Lequeux et al 1995; Heckman & Leitherer 1997; Kunth et al 1998; Gonzalez-Delgado et al 1998a). Moreover, as discussed by Heckman et al (1998), there is a strong correlation between the equivalent widths of these $`UV`$ absorption-lines and the reddening in the $`UV`$ that is analogous to the correlation between reddening in the optical and the $`NaD`$ equivalents widths. The $`IUE`$ spectra discussed by Heckman et al (1998) do not resolve the $`UV`$ absorption-lines, and so can not verify that the correlation is primarily driven by the covering fraction of the absorbing dusty material (as in the case of the $`NaD`$ line). As the archive of $`HST`$ $`UV`$ spectra of starbursts grows, it will be possible to test this.
Images of several edge-on starburst and star-forming galaxies show far-IR and/or sub-mm emission extending one or two kpc along the galaxy minor axis (Alton, Davies, & Bianchi 1999; Alton et al 1998). Multi-color optical images show that kpc-scale extraplanar dust filaments are common in star-forming edge-on galaxies (Howk & Savage 1997, 1999; Sofue, Wakamatsu, & Malin 1994; Phillips 1993; Ichikawa et al 1994). Imaging polarimetry reveals light scattered by extraplanar dust in starburst galaxies (Scarrott, Eaton, & Axon 1991; Scarrott et al. 1993; Scarrott, Draper, & Stockdale 1996; Alton et al 1994; Draper et al 1995). Zaritsky (1994) finds evidence for very extended dust in the halos of spiral galaxies based on the possible detection of reddening in background field galaxies.
As discussed by Howk & Savage (1997) and Aguirre (1999b), there are a variety of mechanisms by which an episode of intense star-formation could lead to the outflow of dust grains. Radiation pressure can “photo-levitate” the grains (Ferrara et al 1991; Ferrara 1998; Davies et al 1997), the Parker instability could help loft material out of the starburst disk (e.g. Kamaya et al 1996), or cold, dusty gas in and around the starburst could be entrained and accelerated outward by the hot outflowing X-ray gas in the superwind (Suchkov et al 1994, and see section 4.2 above).
The superwind mechanism is of the most direct relevance to the present paper, so we briefly evaluate its plausibility. First, we can show that the outward force of the wind on even the largest grains will exceed the inward force of gravity on the grain. For an isothermal galactic potential this force-ratio at a distance $`r`$ from the starburst is given by:
$$F_w/F_g=3\dot{M}v_{term}/4\mathrm{\Omega }_wrv_{rot}^2a\rho $$
(11)
where $`a`$ and $`\rho `$ are the radius and density of the grain (we take $`\rho `$ = 2 gm cm<sup>-3</sup> as representative). For the estimated properties of typical outflow sources in our sample ($`\dot{M}`$ 25 M per year, $`v_{term}`$ 600 km s<sup>-1</sup>, and $`\mathrm{\Omega }_w/4\pi `$ 0.4), $`F_w/F_g`$ will be greater than unity for grains smaller than 7 $`\mu `$$`m`$ ($`r`$/10 kpc)<sup>-1</sup>.
Next, we follow Aguirre (1999b) and estimate the ratio of the sputtering and outflow times for graphite grains immersed in a hot galactic wind ($`\tau _{sp}`$/$`\tau _{out}`$). For the typical parameters we deduce for the outflows in our sample (see above), this ratio is $`\tau _{sp}`$/$`\tau _{out}`$ = 4 ($`a`$/0.1$`\mu `$$`m`$)($`r`$/10 kpc). Thus, large grains could in fact survive the journey to the galactic halo and beyond. The survivability of grains may actually be higher than the above simple estimate if the grains are imbedded inside cold gas clouds propelled by the hot outflow (so that the grains are not directly exposed to the hot gas).
If starburst and star-forming galaxies are indeed capable of ejecting substantial quantities of dust, this could have a profound impact on observational cosmology (e.g. Heisler & Ostriker 1988; Davies et al 1997; Ferrara 1998; Ferrara et al 1999; Aguirre 1999a,b). However, to date, the direct evidence for the existence of intergalactic dust is very sparse. Thermal far-IR emission has been detected from the ICM of the Coma cluster (Stickel et al 1998), and a possible deficit of background QSO’s seen through foreground galaxy clusters has been reported (Romani & Maoz 1992; but see Maoz 1995).
Aguirre (1999a,b) has recently calculated that a dusty inter-galactic medium with $`\mathrm{\Omega }_{dust}`$ = few $`\times `$ 10<sup>-5</sup> would have a visual extinction ($``$ 0.5 magnitudes out to z = 0.7) that would be sufficient to reconcile the Type Ia supernova Hubble diagram (Reiss et al 1998; Perlmutter et al 1999) with a standard $`\mathrm{\Omega }_M`$ = 1, $`\mathrm{\Omega }_\mathrm{\Lambda }`$ = 0 cosmology. Data on the optical colors of high-redshift supernovae show no evidence for reddening, but Aguirre argues that intergalactic dust will have a much greyer extinction curve than standard Galactic dust. This is plausible because small grains will be more easily destroyed by sputtering during and after their journey into the IGM (see above).
In this context, it is instructive to estimate the cosmic mass density of dust grains by the type of outflows investigated in this paper. Aguirre (1999b) has considered this in more detail from a somewhat different perspective, but comes to rather similar conclusions. Let us assume that superwinds associated with the formation of galactic spheroids propelled dust and gas-phase metals into the ICM and IGM, with an amount proportional to the mass in the present-day stars in such systems. We further assume that the mass fractions of the metals locked into grains in the ICM and IGM are $`f_{g,icm}`$ and $`f_{g,igm}`$ respectively. These assumptions imply:
$$\mathrm{\Omega }_{dust,igm}=f_{g,igm}(1f_{g,icm})^1\mathrm{\Omega }_{spheroids}\mathrm{\Omega }_{icm}Z_{icm}\mathrm{\Omega }_{stars,cl}^1$$
(12)
Following FHP and Renzini (1997), we take $`Z_{icm}`$ = 6.7 $`\times `$ 10<sup>-3</sup> (1/3 solar metallicity), $`\mathrm{\Omega }_{spheroids}`$ = 0.0026 $`h_{70}^1`$, $`\mathrm{\Omega }_{icm}`$ = 0.0026 $`h_{70}^{1.5}`$, and $`\mathrm{\Omega }_{stars,cl}`$ = 0.00043 $`h_{70}^1`$. This implies $`\mathrm{\Omega }_{dust,igm}`$ = 1.0$`\times `$ 10<sup>-4</sup> $`f_{g,igm}(1f_{g,icm})^1h_{70}^{1.5}`$. For a normal Galactic dust/metals ratio ($`f_g`$ 0.5), the implied value for $`\mathrm{\Omega }_{dust,igm}`$ is twice as large as the value needed to explain the Type Ia supernova-dimming (Aguirre 1999a). Given the higher densities (and thus, faster grain sputtering times) in the ICM compared to the IGM, we might expect that $`f_{g,icm}<f_{g,igm}`$. More importantly, it is also possible that $`f_{g,igm}<<`$ 0.5 due to the destruction of dust in superwinds and/or the IGM (but see Aguirre 1999b for an optimistic assessment). While the above estimate for $`\mathrm{\Omega }_{dust,igm}`$ should therefore probably be regarded as an absolute upper bound, it is an intriguingly large one from a cosmological perspective.
Finally, we note that since intergalactic dust will emit as well as absorb, its amount is constrained by the cosmic background measured by $`COBE`$ (Ferrara et al 1999). Indeed, Aguirre & Haiman (2000) argue that a significant fraction of the detected cosmic far-IR and sub-mm background must have an intergalactic origin if this dust is abundant enough to strongly affect the Type Ia supernova Hubble Diagram.
### 4.6 Relationship to “Associated Absorption” in AGN
Over the past few years, it has become increasingly clear that a young stellar population is present in the circumnuclear region of a significant fraction of type 2 Seyfert galaxies (e.g. Heckman et al 1995,1997; Gonzalez-Delgado et al 1998b; Schmitt et al 1999; Oliva et al. 1999). Most recently, a near-UV spectroscopic survey of a complete sample of the brightest type-2 Seyfert nuclei by Gonzalez-Delgado, Heckman, & Leitherer (2000) finds direct evidence for hot, young stars in roughly half of the nuclei. In this paper we have established that starbursts drive outflows of cool or warm gas with total column densities of a few $`\times `$ 10<sup>21</sup> cm<sup>-2</sup>, velocities of a few hundred km s<sup>-1</sup>, and covering factors along the line-of-sight of typically 50%. The implication then is that this absorbing material should be detectable in those Seyfert nuclei that also contain a circumnuclear starburst.
In the standard “unified” scenario, type 1 and type 2 Seyfert nuclei are drawn from the same parent population, with the former viewed from a direction near the polar axis of an optically and geometrically-thick “obscuring torus” and the latter from a direction near the equatorial plane of the torus (e.g. Antonucci 1993 and references therein). Thus, in type 1 Seyferts, any starburst-driven outflow could be observed in absorption against the bright nuclear continuum source. While the total column density of the outflowing gas should be similar to the flows studied in this paper, the gas would be exposed to the intense ionizing continuum from the central nucleus, and therefore would be significantly more highly-ionized.
This can be quantified as follows. The ionization state of photoionized gas is determined by the ionization parameter:
$$U=Q/4\pi r^2nc$$
(13)
where $`Q`$ is the production rate of ionizing photons and $`n`$ is the electron density in the photoionized material located a distance $`r`$ from the source. The radial density gradients observed in starburst-driven outflows are consistent with predictions for clouds subjected to the ram pressure associated with the superwind (Heckman, Armus, & Miley 1990):
$$2n(r)kTP(r)=\dot{p}/\mathrm{\Omega }_wr^2$$
(14)
where $`\dot{p}`$ is the rate at which the starburst feeds momentum into the superwind. For photoionized gas, $`T`$ 10<sup>4</sup> K (e.g Osterbrock 1989), so equations 13 and 14 together imply that the magnitude of $`U`$ is set by $`Q/\dot{p}`$, and that $`U`$ will be independent of $`r`$ (neglecting radiative transfer effects). We adopt a generic Leitherer & Heckman (1995) starburst model (Salpeter IMF extending up to 100 M and a starburst lifetime of a few $`\times `$ 10<sup>7</sup> years), include sources of ionization due to both a starburst ($`Q_{}`$) and the type 1 Seyfert nucleus ($`Q_{sy1}`$). For a starburst and type 1 Seyfert nucleus of the same bolometric luminosity, $`Q_{sy1}/Q_{}`$ would be a factor of several. We then obtain the following estimate for $`U`$:
$$U=2.6\times 10^3(\mathrm{\Omega }_w/4\pi )(1+Q_{sy1}/Q_{})$$
(15)
The predicted properties of the absorbing material then overlap significantly with the “associated absorbers” seen in $`UV`$ spectra of type 1 Seyfert nuclei (e.g. Crenshaw et al 1999; Kraemer et al 1999): an incidence rate of roughly 50%, a high line-of-sight covering fraction, outflow velocities of 10<sup>2</sup> to 10<sup>3</sup> km s<sup>-1</sup>, and inferred ionization parameters of $``$ 10<sup>-2</sup>. Crenshaw et al (1999) find an essentially one-to-one correspondence between the presence of $`UV`$ absorption-lines and soft X-ray absorption by hotter and more highly ionized material (the “warm absorber”). We speculate that the hotter and more tenuous phases of the starburst superwind could contribute to the warm absorber.
We emphasize that we are not proposing that all of the “associated absorption” seen in type 1 Seyfert nuclei is produced by gas in a starburst-driven outflow. In some cases, rapid variability or the presence of absorption out of highly-excited lower levels imply densities that are orders-of-magnitude higher than would be tenable for material in a starburst-driven outflow (Crenshaw et al 1999 and references therein). However, it appears that the absorbing material in type 1 Seyfert nuclei can span a broad range in physical and dynamical conditions (Kriss et al 2000). Given important roles for starbursts in the Seyfert phenomenon and for superwinds in the starburst phenomenon, significant absorption due to the superwind material seems unavoidable in some Seyfert nuclei.
## 5 Conclusions
We have discussed the results of moderate-resolution ($`R`$ = a few thousand) spectroscopy of the $`NaI\lambda \lambda `$5890,5896 ($`NaD`$) absorption-line in a sample of 32 far-IR-selected starburst galaxies. These galaxies were selected from either the far-IR-warm sample of Armus, Heckman, & Miley (1989) or the edge-on sample of Lehnert & Heckman (1995), and together span a range from 10<sup>10</sup> to few $`\times `$ 10<sup>12</sup> L in IR luminosity. We found that the stellar contribution to the $`NaD`$ absorption-line is negligible ($`<`$10%) in some objects, but significant ($``$ 70%) in others. We have thus divided our sample into 18 interstellar-dominated (“ISD”) objects ($`<`$ 30% stellar contribution) and 14 strong-stellar-contamination (“SSC”) objects ($`>`$ 40% stellar contribution).
The $`NaD`$ line lies within 70 km s<sup>-1</sup> of $`v_{sys}`$ in all the SSC objects (consistent with a predominantly stellar origin). The $`NaD`$ lines in the SSC nuclei are about 0.2 dex narrower than expected for dynamics of the old stellar population in the bulges of normal galaxies of similar disk rotation speed and Hubble type. Thus, dynamically “cold” material (red supergiants and/or interstellar gas) in the inner part of the starburst makes a significant contribution to the observed $`NaD`$ line in these nuclei.
The kinematics of $`NaD`$ line are markedly different in the ISD objects. The $`NaD`$ line is blueshifted by $`\mathrm{\Delta }`$$`v>`$ 100 km s<sup>-1</sup> relative to the galaxy systemic velocity in 12 of the 18 cases (the “outflow sources”), and the outflow can be mapped over a region of a few-to-ten kpc in size. In contrast, no objects in our sample showed a net redshift in $`NaD`$ of more than 100 km s<sup>-1</sup>. The outflow sources are galaxies systematically viewed more nearly face-on than the other galaxies in our sample: 69% of the galaxies with a ratio of semi-major to semi-minor axes $`a/b`$ 2.0 show $`NaD`$ outflows, while this is true for only 6% of the flatter (more highly inclined) galaxies. This is consistent with the absorbing material being accelerated out along the galaxy minor axis by a bipolar superwind. The absorbing material typically spans the velocity range from near the galaxy systemic velocity ($`v_{sys}`$) to a maximum blueshift of 300 to 700 km s<sup>-1</sup>. We therefore suggest that the outflowing superwind ablates the absorbing gas from ambient clouds at $`v_{sys}`$, and then accelerates it up to a terminal velocity similar to the wind speed. We found no correlation between the widths of the H$`\alpha `$ emission-line and the $`NaD`$ absorption-line subsamples. Evidently, the dynamics of the more tenuous absorbing gas is largely decoupled from that of the dense (high emission-measure) gas that provides most of the nuclear line-emission.
The ratio of the equivalent widths of the two members of the $`NaD`$ doublet ($`R`$) ranges from 1.1 to 1.7 in the ISD sample, implying that the doublet is optically-thick. However, $`R`$ does not correlate with the residual relative intensity at the “bottom” of the stronger $`\lambda `$5890 line profile ($`I_{5890}`$), which ranges from 0.14 (nearly black) to 0.7. Thus, the optically-thick gas does not fully cover the emitting stars (covering factor $``$ 1 - $`I_{5890}`$). The observed equivalent width of the $`NaD`$ line is then set by the product of velocity dispersion and covering factor for the absorbing gas, and we showed that the latter quantity is the dominant one. Using two variants of the classic doublet-ratio technique, we estimated that the $`NaI`$ column densities are $`logN_{NaI}`$ = 13.5 to 14 cm<sup>-2</sup>. This is roughly consistent with column densities measured in a few cases for $`KI`$ using the optically-thin $`\lambda `$$`\lambda `$7665,7699 Å doublet (assuming a solar $`NaI/KI`$ ratio). The total gas columns are uncertain, but the empirical correlation between $`N_{NaI}`$ and $`N_H`$ in the ISM of the Milky Way implies $`N_H`$ few $`\times `$ 10<sup>21</sup> cm<sup>-2</sup>.
We found a strong correlation in the ISD sample between the reddening of the observed stellar continuum and the depth of the $`NaD`$ absorption-line, and a significant but weaker correlation of the line-depth with the reddening of the Balmer emission-lines. Evidently, the gas responsible for the $`NaD`$ absorption is very dusty. The typical implied reddening is $`E(BV)`$ 0.3 to 1 magnitudes over regions several-to-ten kpc in size. For a normal dust-to-gas ratio, the corresponding column densities are $`N_H`$ few $`\times `$ 10<sup>21</sup> cm<sup>-2</sup> (in agreement with the above estimate).
The inferred column densities and measured outflow velocities and sizes imply that the typical mass and kinetic energy associated with the absorbing gas is of-order 10<sup>9</sup> M and 10<sup>56</sup> erg, respectively. The estimated outflow rates of mass and energy are typically 10 to 100 M per year and 10<sup>41</sup> to 10<sup>42</sup> erg s<sup>-1</sup>. The mass outflow rates are comparable to the estimated star-formation rate, and much larger than the rate at which massive stars are returning mass to the ISM. Thus, powerful starbursts can eject as much gas as is being converted into stars, and most of this gas is ambient material that has been “mass-loaded” into the hot gas returned directly by supernovae and stellar winds. The energy outflow rates in the absorption-line gas are of-order 10<sup>-1</sup> of the rate at which massive stars supply mechanical energy. Most of the energy returned by massive stars probably resides in the kinetic and thermal energy of the much hotter X-ray-emitting gas. We showed that the overall properties of the absorbing gas in the outflow sources can be easily reproduced in the context of simple analytic estimates for the properties of interstellar clouds accelerated by the ram pressure of the hot high-speed wind seen via its X-ray emission. Detailed hydrodynamical simulations of galactic winds, while still missing some essential physics, also predict the observed properties of the cool absorbing gas.
We have discussed the implications of our results for the chemical evolution of galaxies and the intergalactic medium. The estimates derived for $`v_{term}`$ using the $`NaD`$ line in the outflow sources agree reasonably well with the outflow speeds implied for an adiabatic wind “fed” by hot gas whose temperature is measured by the observed X-ray-emitting gas. The typical implied values are 300 to 800 km s<sup>-1</sup>, and are independent of the rotation speed of the “host galaxy” over the range $`v_{rot}`$ = 30 to 300 km s<sup>-1</sup>, confirming and extending the result in Martin (1999) based on X-ray data alone. This strongly suggests that the outflows selectively escape the potential wells of the less massive galaxies. We considered a simple model based on Lynden-Bell (1992) in which the fraction of starburst-produced metals that are retained by a galaxy experiencing an outflow is proportional to the galaxy potential-well depth for galaxies with $`v_{esc}<v_{term}`$, and asymtotes to full retention for the most massive galaxies ($`v_{esc}>v_{term}`$). For $`v_{term}`$ in the range we measure, such a simple prescription can reproduce the observed mass-metallicity relation for elliptical galaxies and deposit the required amount of observed metals in the intra-cluster medium. If the ratio of ejected metals to stellar spheroid mass is the same globally as in clusters of galaxies, we predicted that the present-day mass-weighted metallicity of an intergalactic medium with $`\mathrm{\Omega }_{igm}`$ = 0.015 will be $``$ 1/6 solar (see also Renzini 1997).
We have summarized the evidence that starbursts are ejecting significant quantities of dust, emphasizing the results from the present paper. If this dust can survive a trip into the intergalactic medium and remain intact for a Hubble time, we estimated that the upper bound on the global amount of intergalactic dust is $`\mathrm{\Omega }_{dust}`$ $``$ 10<sup>-4</sup>. While this is clearly an upper limit, it is a cosmologically interesting one: Aguirre (1999a,b) argues that dust this abundant could in principle obviate the need for a positive cosmological constant, based on the Type Ia supernova Hubble diagram.
Finally, given the mounting evidence for a connection between starbursts and the Seyfert phenomenon, we have suggested that outflows like those studied here may account for some (but not all) aspects of the “associated absorption” seen in type 1 Seyfert nuclei.
We would like to thank Ken Sembach for useful on-going discussions and advice. Discussions with David Neufeld, Mark Voit, Don York, and Donna Womble were helpful during the formative stages of the project. The partial support of this project by NASA grant NAGW-3138 is acknowledged.
Note. Col. (2) — Galaxy systemic velocity (km s<sup>-1</sup>) in the heliocentric frame. In order of preference, these are determined from: galaxy rotation curves (r), global $`CO`$ 115 GHz emission-line profiles (c), nuclear stellar velocities (s), global $`HI\lambda `$21cm emission-line profiles (h), and nuclear optical emission-line profiles (e). The rotation curve velocities (r) are taken from LH95 except for NGC2146 from Prada et al (1994). Items marked ‘n’ come from NED. Items marked ‘e’ or ‘s’ are based on data obtained during the observing runs discussed in the present paper. Items marked ‘c’ are: NGC 660 (Elfhag et al 1996; Young et al 1995), NGC 1614 (Elfhag et al 1996; Young et al 1995; Aalto et al 1991; Sanders, Scoville, & Soifer 1991; Casoli et al 1991), M 82 (Lo et al 1987), IRAS10173+0828 (Planesas, Mirabel, & Sanders 1991), NGC 3256 (Aalto et al 1991; Casoli et al 1991; Mirabel et al 1990), IRAS10565+2448 (Downes & Solomon 1998), and Arp 220 (Young et al 1995; Solomon, Downes, & Radford 1992). Based on the intercomparison of independent measurements for a given galaxy, the typical uncertainties in $`v_{sys}`$ range from 10 km s<sup>-1</sup> for the nearby, relatively normal galaxies to as much as 100 km s<sup>-1</sup> for the most distant systems (generally, highly disturbed mergers). Col. (3) — Total infrared luminosity from 8 to 100 microns, based on IRAS data and the definition of L<sub>IR</sub> given in Sanders & Mirabel (1996). We assume throughout that $`H_0`$ = 70 km s<sup>-1</sup> Mpc<sup>-1</sup>. Col. (4) — Blue absolute magnitude for the galaxy, corrected for foreground (Galactic) extinction, but not for internal extinction. Taken from LH95 when available (adjusted to $`H_0`$ = 70 km s<sup>-1</sup> Mpc<sup>-1</sup>) or based on the data in NED or Armus, Heckman, & Miley (1987). Col. (5) — The ratio of the optical semi-major to semi-minor axes. These are taken (in order of preference) from LH95, the images published in Armus, Heckman, & Miley (1987; 1990), or NED. For the highly disturbed merging systems, we have measured this ratio at intermediate radii (excluding both faint tidal tails and the inner regions where dust obscuration is most significant). Col. (6) — The amplitude of the rotation speed of the galaxy. In order of preference, we have based these on rotation curves (“r” from LH95), global $`HI\lambda `$21cm profiles corrected for inclination and turbulence (“h” - see LH95 for details), and global $`CO`$ 115 GHz emission-line profiles using the half-width at 20% of the peak intensity and then correcting for inclination (“c” - using the same data as in Column 2). For M 82 we have replaced the value listed in LH95 by the more recent determination by Sofue (1998). The uncertainties in most cases are dominated by the inclination correction. We estimate the resulting uncertainties to be $`<`$ 0.1 dex for all but the cases of mergers and strongly interacting galaxies, where the inclination corrections lead to an uncertainty of roughly 0.2 dex (denoted by :). Col. (7) — Sample from which the galaxy was drawn (Armus, Heckman, & Miley 1989; Lehnert & Heckman 1995). Col. (8) — Observing runs used in this paper (see Table 2).
Note. Col. (2) — The estimated contribution to the observed $`NaD`$ line by cool stars (the remainder is interstellar in origin). See text for details. Galaxies with $`f_{}`$ 40% are members of the strong-stellar-contamination sample (SSC), while those with $`f_{}`$ 30% are members of the interstellar-dominated (ISD) sample. Based on the agreement between $`f_{}`$ determined by the two independent techniques discussed in the text, the uncertainty is typically $`\pm `$10% . Col. (3) — Heliocentric velocity of the $`NaD`$ absorption-line (km s<sup>-1</sup>) measured by fitting the profile with a pair of Gaussians constrained to have the separation in wavelength appropriate for the red-shifted $`NaD`$ doublet. Based on comparison of independent measurements of this quantity in the cases for which we have multiple spectra, we estimate that the typical measurement uncertainty is $`\pm `$ 20 km s<sup>-1</sup>. Col. (4) — The velocity difference between the $`NaD`$ absorption-line and the galaxy systemic velocity in the galaxy rest-frame: $`\mathrm{\Delta }`$v = (v<sub>NaD</sub> \- v<sub>sys</sub>)/(1 + v<sub>sys</sub>/c). The relevant quantities are given in Col.2 of Table 1 and Col. 3 of this Table. A typical uncertainty in this velocity difference is $`\pm `$ 20 km s<sup>-1</sup>, for the relatively bright nearby galaxies with well-determined values for $`v_{sys}`$ up to $`\pm `$ 100 km s<sup>-1</sup> for the most distant and far-IR-luminous galaxies (highly disturbed mergers with uncertain $`v_{sys}`$). Col. (5) — The full-width at half-maximum of each of the two Gaussians fit to the doublet (km s<sup>-1</sup>). $`W`$ was constrained to be the same for the two doublet members. The listed value has had the instrumental contribution to the measured value removed by assuming the intrinsic and instrumental widths add in quadrature: W = \[W$`{}_{}{}^{2}{}_{obs}{}^{}`$ \- W$`{}_{}{}^{2}{}_{instr}{}^{}`$\]<sup>1/2</sup>. Based on comparison of independent measurements of this quantity in the cases for which we have multiple spectra, we estimate that the typical measurement uncertainty is $`\pm `$ 20 km s<sup>-1</sup>. Col. (6) — The rest-frame equivalent width (Å) for the $`NaD`$ doublet. Based on comparison of independent measurements of this quantity in the cases for which we have multiple spectra, we estimate that the typical measurement uncertainty is $`\pm `$ 0.2 Å. Col. (7) — The normalized residual intensity at the center of the $`NaD`$ $`\lambda `$5890 line profile (I<sub>5890</sub> = 0 corresponds to a totally black line center). This has been corrected for the effect of the spectral resolution assuming Gaussian profiles: (1-I<sub>5890</sub>) = (W<sub>obs</sub>/W)(1 - I<sub>5890,obs</sub>). Based on comparison of independent measurements of this quantity in the cases for which we have multiple spectra, we estimate that the typical measurement uncertainty is $`\pm `$ 0.02. Col. (8) — The ratio of the equivalent widths of the $`NaD\lambda \lambda `$5890,5896 transitions. A ratio $`R`$ = 2 (1) corresponds to an optical depth of 0 (infinity). Based on comparison of independent measurements of this quantity in the cases for which we have multiple spectra, we estimate that the typical measurement uncertainty is $`\pm `$ 0.1. Col. (9) — The full-width at half-maximum of a Gaussian fit to the nuclear $`H\alpha `$ emission-line. These are taken from AHM, LH95, or our own unpublished spectra. The listed value has had the instrumental contribution to the measured value removed by assuming the intrinsic and instrumental widths add in quadrature. Typical uncertainties are $`\pm `$ 20 km s<sup>-1</sup>.
Note. Col. (2) — The ratio of the nuclear H$`\alpha `$ and H$`\beta `$ emission-line fluxes. These have been corrected for the effects of underlying stellar absorption-lines (assuming a stellar equivalent width of 1.5 Å) and for foreground Galactic extinction (using a standard Galactic extinction curve and the measured Galactic $`HI`$ column density. The data come from AHM, Veilluex et al (1995), data from runs 3 or 6 (Table 2), Vaceli et al (1997), or Dahari & DeRobertis (1988). Typical uncertainties are $`\pm `$5%. Unreddened ionized gas would have a flux ratio of 2.86 for standard Case B conditions. Col. (3) — The ratio of the flux densities (F<sub>λ</sub>) in the nuclear continuum near the wavelengths of H$`\alpha `$ and H$`\beta `$. The values have been corrected for foreground Galactic extinction (see above). The data come from AHM, Veilluex et al (1995), or our runs 3 or 6 (Table 2). Typical uncertainties are $`\pm `$5%. Note that an unreddened starburst corresponding to constant star-formation for 30 Myr would have an intrinsic color in these units of 0.5. Col. (4) — The projected size (in kpc) of the region along the spectrograph slit exhibiting strongly blueshifted (by $`>`$ 100 km/s) $`NaD`$ absorption. In NGC1572, NGC1614, NGC3256, NGC7552, and NGC7582 we have measured this along two position angles. Col. (5) — An estimate of the terminal velocity implied by the $`NaD`$ absorption-line profile (v<sub>term</sub> = $`\mathrm{\Delta }`$ v + 0.5 W). See Table 3. Col. (6) — The rotation speed of the starburst galaxy. See Table 1. |
warning/0002/hep-ex0002032.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Several theories and many models predict the existence of magnetic monopoles (MMs) \[1-5\] with a magnetic charge which obeys the Dirac quantization rule , $`eg=n\mathrm{}c/2`$, where $`n`$ is an integer, $`n`$=1,2,3… For $`n`$=1, assuming that the basic electric charge is that of the electron, one has $`g=g_D=\mathrm{}c/2e=3.2910^8`$ u.e.s. Most theoretical papers concern superheavy magnetic monopoles predicted by Grand Unified Theories (GUT) of electroweak and strong interactions; there are also many papers concerning still heavier MMs, and others concerning MMs with intermediate and low masses ; if the basic electric charge is that of the quark $`d`$ the basic magnetic charge increases by a factor of three.
A MM and an atomic nucleus N may form a bound system. Some models predict the existence of states with both electric and magnetic charges (dyons) . As far as the energy losses in matter, a dyon and a $`(MM+N)_{bound}`$ system behave in the same way. The system $`(MM+N)_{bound}`$ can be produced mainly via radiative capture , $`MM+N(MM+N)_{bound}+\gamma `$, and the bound system may be subject to photodissociation , $`(MM+p)_{bound}+\gamma MM+p`$. Since the cross section for radiative capture is $`1<\sigma <10mb`$ for MMs with $`10^4<\beta <10^3`$, it is possible that many MMs in the cosmic radiation can be in the state $`(MM+p)_{bound}`$. Neglecting the catalysis of proton decay, the MMs or the $`(MM+p)_{bound}`$ systems with $`\beta 10^3`$ traversing rock, probably reach an equilibrium mixture of about $`50\%`$ MMs and $`50\%`$ $`(MM+p)_{bound}`$ . The minimum velocity of the $`(MM+p)_{bound}`$ system for which one may have a break-up reaction is $`\beta 10^3`$ ( at this value the three particles in the final state (MM, proton and target nucleus) are essentially at rest in the center of mass frame of the MM).
Most experimental searches concern superheavy GUT magnetic monopoles and most of the published flux limits apply to an isotropic flux of bare monopoles of unit magnetic charge (the monopole catalyzed nucleon decay is not considered). In order to search for MMs and dyons over a large $`\beta `$-range it is necessary to know how these particles loose energy in different types of detectors and if they can cross the Earth . There has been increasing interest in lower mass monopoles, such as those discussed in Refs. . Proposals have been made to search for them at sea level and at high altitude and also at the future LHC collider at CERN .
This lecture is an update of the lectures on MMs at the 1995 Trieste 4<sup>th</sup> School on Non Accelerator Particle Astrophysics . Here there will be a more complete analysis of the energy losses of MMs and of dyons in different detectors and in the Earth; an update is made of the GUT MM searches. An extensive bibliography on MMs is given in Ref. . As a byproduct we shall discuss the searches and the present limits for nuclearites. These limits are also valid for charged Q-balls, aggregates of supersymmetric particles .
## 2 Energy losses of monopoles and dyons in detectors
One has to consider the total energy losses of MMs and of dyons in the detectors and also that fraction of the energy loss which leads to detection. In streamer tubes this fraction is the ionization energy loss in the gas, in scintillators it is the excitation energy loss which leads to the emission of light, in nuclear track detectors it is the restricted energy loss, t. i. the energy loss contained in a 10 nm diameter around the MM trajectory. In Ref. a thorough analysis is made of these losses and of their dependence on the MM velocity v=$`\beta `$c.
### 2.1 Total energy losses in detectors
Let us consider first the scintillators and the streamer tubes. At high velocity ($`\beta >0.05`$) the total energy losses of MMs in scintillators and streamer tubes are mainly due to ionization; they may be calculated using the Bethe-Bloch formula. At intermediate $`\beta `$ (few $`10^3<\beta <10^2`$) the main contribution to energy losses is again due to ionization, and one uses the formula of Ahlen-Kinoshita for nonconductors, corrected by Ritson :
$$\left(\frac{dE}{dx}\right)_{\text{ionization}}2.610^6\left(\frac{g}{g_D}\right)^2\beta ^{1.7}\frac{\text{MeV}\text{g}}{\text{cm}^2}$$
(1)
The results are shown as curves A in Fig. 1a for MMs in scintillators.
At low velocity ($`10^5<\beta <10^3`$) the dominant energy loss mechanism is that of elastic collisions of the MM with the atoms. This contribution may be computed by numeric procedures, neglecting the hydrogen atoms (which means that one obtains a lower limit for the energy losses). The results are shown as curves B in Fig. 1a for MMs in scintillators .
### 2.2 Light yields in scintillators for MMs
For MMs the light yield $`dL/dx`$ in a scintillator is related to the energy loss $`dE/dx`$ by
$$\frac{dL}{dx}=A\left[\frac{1F}{1+AB(1F)(dE/dx)}+F\right]\frac{dE}{dx}$$
(2)
where $`dE/dx`$ is the total electronic energy loss. For relatively small $`dE/dx`$ the light yield is proportional to the energy loss with the proportionality constant $`A`$. $`B`$ is the the so-called quenching parameter: the light yield from the energy deposited near the track (the first term in the parenthesis of Eq. 2) saturates for large energy losses. For $`\beta >0.1`$ some electrons have sufficient energy to escape from the region near the track core ($`\delta `$ rays); $`F`$ is the fraction of energy loss which results from excitations outside the core; these excitations are assumed not to be quenched.
The energy loss $`dE/dx`$ for a monopole is given at high velocity ($`\beta >0.05`$)) by the corrected Bethe-Bloch formula.
At lower velocities ($`\beta <0.01`$) the energy losses of a MM in scintillators are $`\frac{1}{4}(g/g_D)^2`$ times the energy losses of a proton of the same $`\beta `$; these are computed from the energy losses in hydrogen and carbon, adding an exponential factor from the fit to the low-velocity proton data. For $`0.01<\beta <0.05`$ a smooth interpolation is used. The results of the calculations are shown in Fig. 1b. Notice that at very low velocities ($`0.0002<\beta <0.0005`$ for $`g=g_D`$) the light yield increases with $`\beta `$, then it saturates at a value of 1.2 MeV/cm (region A of Fig. 1b). The increase in the light yield observed in region B of Fig. 1b, is due to changes in the quenching parameters. For $`\beta 0.09`$ there is production of delta rays, and the light yield increases again with $`\beta `$ (region C of Fig. 1b).
For $`0.003<\beta <0.1`$ the light yield is thus independent of the magnetic charge value; at low and high $`\beta `$ the light yield increases quadratically with $`g`$. The light yields in NE110 and in the MACRO liquid scintillator are equal within a few percent and only one curve is drawn. The light yield at low beta is a lower limit, since the present calculations do not take into account possible contributions that arise from the mixing and crossing (Drell effect) of molecular electronic energy levels at the passage of the magnetic charge; this could result in molecular excitations and emission of light.
The Ahlen-Tarlé curve shown in Fig. 1b, refers to the calculation of Ref. which predicts a kinematical cut off at $`\beta =710^4`$ due to the excitation energy of benzene molecules (5 eV); however experiments with low energy protons have shown that no such cut off occurs. The light yield is a small fraction of the total energy loss, see Figs. 1a and 1b.
### 2.3 Energy losses of MMs in streamer tubes
At high velocities, $`\beta >0.05`$, the Bethe Bloch formula holds for the energy losses of MMs in the streamer tubes. The resulting energy losses are given in Fig. 2a, as curves A.
At low velocities, $`10^4<\beta <10^3`$, and for $`g=g_D`$ in helium, the Drell effect occurs: the energy levels of the helium atoms are changed by the presence of a magnetic charge, and at the passage of the magnetic charge an electron can make a transition to an excited level:
$$\left(\frac{dE}{dx}\right)_{Drell}=11(\beta /10^4)[1(9.310^5/\beta )^2]^{3/2}\text{MeV cm}^2\text{/g}$$
(3)
Such energy losses lead to the atomic excitation of helium atoms; these lead in turn to the ionization of the n-pentane molecules via the Penning effect.
At intermediate velocities, $`210^3<\beta <710^2`$, one considers the medium as a degenerate electron gas. The results of the calculations are shown in Fig. 2a, as curves B.
For dyons one has to sum up the direct excitation and ionization produced by the electric charge and by the moving magnetic charge. In addition the elastically recoiling hydrogen and carbon nuclei add contributions. The results of the calculation for the energy losses of dyons are shown in Fig. 2b.
### 2.4 Energy losses of MMs and of dyons in the CR39 nuclear track detector
The restricted energy losses of MMs and of dyons as a function of $`\beta `$ in the nuclear track detetctor CR39 have been computed using similar approximations and are presented in Figs. 3a and 3b.
## 3 Energy losses of MMs and of Dyons in the Earth.
In order to compute the acceptance of a detector to MMs coming from above and from below one needs a model of the Earth mantle and nucleus, and an evaluation of the energy losses down to very small velocities. The computations of Ref. show that only heavy monopoles can traverse the Earth: for istance for $`g=g_D`$ and $`\beta `$ = $`10^3`$ only MMs with $`m_M>10^{14}`$ GeV can traverse the Earth. Fig. 4 shows the accessible mass region for MMs of different velocities coming from above; for $`g=g_D`$ and $`\beta =10^3`$ a MM must have $`m_M>10^{10}`$ GeV, $`10^6`$ GeV, $`10^5`$ GeV in order to reach the underground MACRO detector, the Earth surface, and a detector at 5230 m altitude, respectively .
## 4 Experimental searches for supermassive GUT monopoles
A flux of cosmic GUT supermassive magnetic monopoles may reach the Earth and may have done so for the whole life of the Earth. The velocity spectrum of these MMs could be in the range $`310^5<\beta <0.1`$. Searches for such MMs have been performed with superconducting induction devices whose combined limit is at the level of $`2`$$`\mathrm{\hspace{0.17em}10}^{14}\mathrm{cm}^2\mathrm{s}^1\mathrm{sr}^1`$, independent of $`\beta `$, see Fig. 5. Direct searches were performed underground using scintillators, gaseous detectors and nuclear track detectors (mainly CR39) . No monopoles have been detected, and the present 90% C.L. flux limits are given in Fig. 5. Indirect searches were performed with old mica samples .
## 5 Intermediate mass magnetic monopoles
Magnetic monopoles with intermediate masses, $`10^3<m_M<10^{12}`$ GeV, have been proposed by several authors, see f.e. Ref. . These monopoles could be present in the cosmic radiation, but would not reach underground detectors, see Fig. 4 . Detectors at the Earth surface are capable to detect MMs coming from above if they have masses larger than $`10^510^6`$ GeV; lower mass monopoles may be searched for with detectors located at high mountain altitudes, or even higher, in balloons and in satellites. Few experimental results are available . An experiment has been proposed to search for intermediate mass monopoles with nuclear track detectors at the Earth surface and at high altitude .
## 6 Magnetic monopole searches at accelerators
Classical MMs have been searched for at every new accelerator up to masses of about 1 TeV . The main motivation is connected with the original proposal of Dirac, in 1931 . With the advent of higher energy accelerators, in particular of the LHC, it should be possible to search for MMs with masses up to about 7 TeV, a mass region which may explore the possibility of magnetic monopoles connected with the electroweak unification scale . A simple experiment using nuclear track detectors is proposed to be used with the LHC collider .
Ginzburg and Schiller have proposed to analyze $`\gamma \gamma `$ collisions at high energy colliders, below monopole production threshold . $`\gamma \gamma `$ scattering should be enhanced due to the large coupling of monopoles to photons. At the Tevatron the effect could be seen as pair production of photons with energies of 200-400 GeV and roughly compensated transverse momenta of 100-400 GeV/c. This search could explore monopole masses of 1-2 TeV; at LHC the explored masses could be 7-19 TeV.
## 7 Experimental searches for nuclearites
Nuclearites (strange quark matter) should be aggregates of $`u`$, $`d`$ and $`s`$ quarks in equal proportions and of electrons to ensure their electrical neutrality. They should be stable for all barion numbers in the range between ordinary nuclei and neutron stars ($`A10^{57}`$). They could be the ground state of QCD and could have been produced in the primordial Universe or in violent astrophysical processes. Nuclearites could contribute to the dark matter in the Universe. An upper limit on the nuclearite flux may be estimated assuming that $`\mathrm{\Phi }_{max.}=\rho _{DM}v/(2\pi M)`$, where $`\rho _{DM}10^{24}`$ $`\mathrm{gcm}^3`$ is the local dark matter density; $`M`$ and $`v`$ are the mass and the velocity of nuclearites, respectively .
The most relevant direct upper flux limits for nuclearites come from three large area experiments: the first two use the CR39 nuclear track detector; one experiment was performed at mountain altitude , the second at a depth of 10<sup>4</sup> $`\mathrm{g}\mathrm{cm}^2`$ in the Ohya mines ; the third experiment is the MACRO detector which uses also liquid scintillators besides nuclear track detectors . Indirect experiments using old mica samples could yield the lowest flux limits, but they are affected by inherent uncertainties . Some exotic cosmic ray events were interpreted as due to incident nuclearites, for example the “Centauro” events, the anomalous massive particles, etc. . The interpretation of those possible signals are not unique and the used detectors are not redundant.
The main energy loss mechanism for nuclearites passing through matter is that of atomic collisions. While traversing a medium the nuclearites should displace the matter in their path by elastic or quasi-elastic collisions with the ambient atoms . The energy loss rate is
$$dE/dx=\sigma \rho v^2,$$
(4)
where $`\sigma `$ is the nuclearite cross section, $`v`$ its velocity and $`\rho `$ the mass density of the traversed medium. For nuclearites with masses $`M1.5`$ $`\mathrm{ng}`$ the cross section may be approximated as:
$$\sigma \pi \left(\frac{3M}{4\pi \rho _N}\right)^{2/3}$$
(5)
It is assumed that the density of strange quark matter is $`\rho _N3.510^{14}`$ $`\mathrm{g}\mathrm{cm}^3`$ somewhat larger than that of atomic nuclei. For typical galactic velocities, nuclearites with masses larger than 0.1 g could traverse the Earth. For nuclearites of smaller masses the collisions are governed by their electronic clouds, yielding $`\sigma \pi 10^{16}`$ $`\mathrm{cm}^2`$. Most nuclearite searches were obtained as byproducts of superheavy magnetic monopole searches. The cosmic ray flux limits are therefore similar to those obtained for MMs.
In Fig. 6 is presented a compilation of limits for a flux of downgoing nuclearites compared with the dark matter limit, assuming a velocity at ground level $`\beta =v/c=210^3`$. This speed corresponds to nuclearites of galactic or extragalactic origin. In the figure we extended the MACRO limit above the dark matter bound, in order to show the transition to an isotropic flux for nuclearite masses larger than 0.1 g ($`10^{23}`$ GeV ).
Figure 6: 90% C.L. flux upper limits versus mass for nuclearites with $`\beta =210^3`$ at ground level. These nuclearites could have galactic or extragalactic origin. The limits are from MACRO , from Refs. (“Nakamura”), (“Orito”) and the indirect Mica limits of Ref. .
## 8 Conclusions
Comparing this lecture at the 5<sup>th</sup> School to that at the 4<sup>th</sup> School one notices the improvements made in this field during a 3-year period. The changes concern some theoretical aspects , computation of the energy losses in detectors and in the Earth , new searches for MMs in the cosmic radiation, see Fig. 5 , new methods to search for classical monopoles , etc. . In the fields related to that of MMs, because of similarity of detection techniques, one has theoretical extensions to Q-balls and new limits on nuclearites, see Fig. 6 .
Acknowledgements. We would like to acknowledge the cooperation of many colleagues of the MACRO experiment, in particular M. Giorgini, T. Lari, G. Mandrioli, M. Ouchrif, V. Popa, P. Serra, M. Spurio, and others. |
warning/0002/astro-ph0002063.html | ar5iv | text | # STRUCTURE OF THE CIRCUMNUCLEAR REGION OF SEYFERT 2 GALAXIES REVEALED BY RXTE HARD X–RAY OBSERVATIONS OF NGC 4945
## 1 Introduction
Our current best picture of nuclei of Seyfert galaxies includes the central source (i.e. black hole, accretion disk and broad line region) embedded within an optically thick molecular torus (cf. Antonucci & Miller 1985). The object is classified as a Seyfert 1 for viewing directions which lie within the opening angle of the torus, so that there is a direct view of the nucleus, and as a Seyfert 2 for directions intersecting the obscuring material. The torus absorbs the optical, UV and soft X–ray nuclear light, so the nucleus in Seyfert 2s can only be seen at these energies through scattered radiation. One of the central questions still under debate in the context of the Unification Models of the two types of Seyferts is: how optically thick is this putative torus?
NGC 4945 is a nearby (3.7 Mpc; Mauersberger et al. 1996) edge–on galaxy. It has strong starburst activity, producing intense IR emission concentrated in a compact nuclear region (Rice et al. 1988; Brock et al. 1988), and a “superwind” outflowing along the minor axis of the galaxy (Heckman, Armus, & Miley 1990). It also has an active nucleus, first seen unambiguously in Ginga X–ray observations (Iwasawa et al. 1993), confirming the Seyfert 2-type classification. These data showed a heavily obscured, strong hard X–ray source above 10 keV, confirmed by the CGRO OSSE observations (Done, Madejski, & Smith 1996) which in turn revealed that NGC 4945 is one of the brightest extragalactic sources in the sky at 100 keV! The absorbing column, a few $`\times `$ $`10^{24}`$ cm<sup>-2</sup>, is among the largest which still allows a direct view of the nucleus at hard X–ray energies. If such a scattering medium were to subtend a large solid angle from the nucleus, it would smear out any intrinsic hard X–ray variability on timescales shorter than the light travel time through it.
The presence of the Seyfert nucleus is further supported by the fact that the object is a megamaser source (detected in the H<sub>2</sub>O bands) implying an edge–on geometry, but one of the key features which makes its study so important is that it is one of only four AGN where the black hole mass can be constrained (at $`1.4\times 10^6`$ $`M_{}`$) from detailed mapping of the megamaser spots (Greenhill, Moran, & Herrnstein 1997). As such, it is one of a few unique sources where the luminosity in Eddington units can be reliably estimated.
Below we present the data from the RXTE, confirming the large absorbing column, but also revealing large amplitude hard X–ray variability on a time scale of days. A distant absorber with an appreciable optical depth, subtending a large solid angle as seen by the nucleus, would smear out the rapid variability on time scales shorter than the light travel time through such an absorber. Given our data, we conclude that the optically thick absorber cannot be both distant and geometrically thick.
## 2 Observations: Spectrum and Variability
NGC 4945 was observed by the Rossi X–ray Timing Explorer (RXTE) satellite for about a month, starting on October 8, 1997. The observations included 38 pointings of $`2000`$ s each, taken about once per day. The Proportional Counter Array (PCA) and High Energy X–ray Timing Experiment (HEXTE) data were reduced using standard procedures. For the PCA data, this included an extraction in the standard2 mode using the ftool saextrct; the estimation of the background was done via ftool pcabackest, using the background model “L7” (model files pca\_bkgd\_faint240\_e03v03.mdl and pca\_bkgd\_faintl7\_e03v03.mdl). For consistency, we use data only from 3 PCA detectors, which were active during all pointings. For HEXTE, the background is collected simultaneously by switching two halves of the array on- and off-source every 16 seconds; both source and background files were extracted using the ftool seextrct, and dead-time corrected using ftool hxtdeadpha. The total “good data” intervals were: 69,280 s for PCA, 18,364 s and 19,060 s for HEXTE clusters A and B.
The summed background-subtracted PCA and HEXTE files were fitted with a phenomenological model including a hard power law with low-energy photoelectric cutoff (using the cross-sections and abundances as given in Morrison & McCammon 1983) and a high-energy exponential cutoff (assumed to be at an energy $`E_c`$ of 100 keV, in agreement with the high energy Seyfert spectra; e.g., Zdziarski et al. 1995). Our model also includes a Gaussian Fe K emission line, plus a soft component, which we modeled as another power law. In our fits, we used the PCA data corresponding to the energy range of 3 to 30 keV, and HEXTE data for 20 to 100 keV. The resulting fit (cf. Fig. 1) was essentially consistent with the Ginga / OSSE results of Done et al. (1996). The hard power law (with $`E_c`$ of 100 keV) showed an energy index of $`0.45\pm 0.1`$, with absorption of $`4.5\pm 0.4\times 10^{24}`$ cm<sup>-2</sup>, and an observed 10 - 50 keV flux of $`1\times 10^{10}`$ erg cm<sup>-2</sup> s<sup>-1</sup>, with resulting $`\chi ^2=80.6/75`$ d.o.f. The Fe K line energy was at $`6.38\pm 0.05`$ keV, with an intrinsic width $`\sigma `$ of $`0.37\pm 0.13`$ keV, and a flux of $`0.9\times 10^4`$ photons cm<sup>-2</sup> s<sup>-1</sup>. Allowing the cutoff to be unconstrained yielded the best fit of $`90_{30}^{+130}`$ keV. Regarding the soft component, its energy power law index of was $`0.57\pm 0.15`$, with the 1 keV monochromatic flux of of 0.001 photons cm<sup>-2</sup> s<sup>-1</sup> keV<sup>-1</sup>. It is important to note that the PCA field of view is about 1 deg<sup>2</sup>, so at least a fraction of the Fe K line and soft component flux could have arisen from the more extended (non-nuclear) region, a likely possibility given the starburst nature of the galaxy. Furthermore, given its modest flux, which over the 2 – 10 keV band is $`6\times 10^{12}`$ erg cm<sup>-2</sup> s<sup>-1</sup> – only three times that of the $`1\sigma `$ fluctuations of the Cosmic X–ray Background on the angular scale of the PCA field of view – we caution that any detailed spectral analysis of the PCA data for this soft component is unreliable. Nonetheless, we can clearly reject the hypothesis that the entire flux of this soft component is due to some kind of a “leaky absorber,” as it does not appear to vary.
The spectral analysis above confirmed the results of Done et al. (1996) that the source spectrum consists of the soft, relatively faint, unabsorbed component, a bright, heavily absorbed (hard) component, and a strong Fe K line. With that, we studied the variability of each component separately. The flux of the soft continuum component (below 8 keV) is consistent with being constant; this is also true of the Fe K line. The hard component (8 – 30 keV), on the other hand, is highly variable, with a factor of 4 change in 10 days, and a factor of 1.7 – 2 in 1 day between the minimum and maximum flux. This is plotted in Fig. 2. This light curve (collected with 3 PCUs, and binned on 1 day intervals) shows RMS variance ($`1\sigma `$) of 0.82 cts s<sup>-1</sup>. Unfortunately, the source was too faint to study the variability with HEXTE.
We investigated if the variability could be due to instrumental effects, and specifically, imprecise background subtraction. To assess this, we also analyzed in an analogous manner the hard (8 – 30 keV) light curves binned in 1 day intervals of a cluster of galaxies Abell 754 (cf. Valinia et al. 1999) and a faint quasar PG1211+143, which shows very little flux in the PCA data above 8 keV (Netzer, Madejski, & Kaspi in prep.). Analysis of 9 data points collected over 9 days for A754 (from which no source variability is expected) yields $`\sigma =0.018`$ cts s<sup>-1</sup>. For PG1211+143, we had 32 pointings spread nearly uniformly over 6 months. These data, where some intrinsic source variability may be present, yield $`\sigma =0.13`$ cts s<sup>-1</sup>; we consider this a conservative upper limit to the instrumental effects, and thus deem the rapid hard X–ray variability of NGC 4945 with $`\sigma `$ of 0.82 cts s<sup>-1</sup> highly significant.
Could this variability be due to varying absorption? We examined this possibility by modeling separate spectra from high and low count rate observations. To improve statistics, we co-added the PCA spectra from a number of individual observations with highest and lowest count rates. The two resulting spectra were then modeled using the Monte Carlo absorption model discussed below. We assumed first that the intrinsic source spectrum (and normalization) is the same for both, but the absorption is different and, secondly, that the normalization of the intrinsic source has changed while the absorption stayed constant. The first hypothesis yielded $`\chi ^2=401/110`$ dof, while the second one yielded $`\chi ^2=143/110`$ d.o.f. This clearly shows that the variability is intrinsic to the unabsorbed nucleus.
## 3 Discussion
The variability we see in NGC 4945 is then entirely compatible with that expected from the intrinsic source, with no significant scattered delay by a distant material. If this absorber has an appreciable Compton thickness (as is the case here), and if it subtends a large solid angle to the X–ray source, then it should intercept and scatter a large fraction of the flux. If it is also distant from the nucleus, the light travel time effects will “wash out” any intrinsic variability on time scales shorter than the light travel time through the absorber. Conversely, a structure with much smaller scale height subtends a much smaller solid angle, making scattering less important. This also precludes the observed X–rays to be purely due to Compton reflection, as this would require a contrived geometry of the reflector with respect to the primary X–ray source: the reflector would have to be located very close to a completely covered central source. We note that by comparison, a well studied unabsorbed Seyfert 1 NGC 3516 showed X–ray variability with a similar fractional amplitude on $`10\times `$ longer time scales (Nandra & Edelson 1999) than seen in NGC 4945. While the black hole mass in NGC 3516 is not as well known, it is estimated to be $`10^7`$ $`M_{}`$. The fact that the ratio of variability time scales is roughly the same as the ratio of nuclear masses – as expected for accreting black holes – further supports our conclusion that the hard X–ray variability we see in NGC 4945 is intrinsic.
With the optical depth to electron scattering of a few, the shape of the absorption cutoff would be different than expected from pure photoelectric absorption, and the detailed shape of the emergent spectrum depends on the geometry. We model this numerically with a Monte Carlo code as given in Krolik, Madau, & Życki (1994), assuming a torus with square cross-section where the half-angle subtended by the torus, $`\theta _0`$, its optical depth to electron scattering $`\tau _\mathrm{e}`$, and the power law index of the incident energy spectrum $`\alpha `$ are free parameters (the Comptonization cutoff is set to 100 keV as above). The results of our fits (using the PCA data over the range of 3 – 20 keV, and HEXTE data as above) are shown in Table 1, where the 90% confidence regions on $`\mathrm{\Gamma }`$ and $`\tau _\mathrm{e}`$ are typically 0.1.
A small scale height absorber ($`\theta _010^{}`$) gives $`\tau _\mathrm{e}=2.4`$, compared to a large scale height ($`\theta _080^{}`$) which gives $`\tau _\mathrm{e}=2.1`$. With an iron abundance of twice Solar, these fits change to $`\tau _\mathrm{e}=1.7`$, and $`\tau _\mathrm{e}=1.5`$, respectively. (Since the photoelectric cutoff present in our data is mainly sensitive to the column density of Fe, larger-than-Solar abundance of Fe would make us overestimate the true absorbing column if we assumed Solar abundances, and vice-versa.) While statistically these might marginally favor the large scale height absorber, we consider that all the fits are probably equally likely given that modeling the spectrum with a fixed cutoff energy may introduce systematic uncertainties. Our calculations include the Fe K emission line produced by the torus but we also include an additional Fe line (such as may be expected to arise in the photoionized scattering medium). Those calculations also imply that large ($`>4\times `$ Solar) Fe abundances can be excluded: they would imply a stronger Fe K line than is seen in our data. In reality, this limit is probably more stringent, given the fact that at least a fraction of the Fe K line originates in a more extended region.
These Monte Carlo results also give the distribution of the number of scatterings which the photons undergo before reaching the observer positioned in the equatorial plane. This is key in determining the solid angle subtended by the optically thick absorber, and thus its vertical size scale. Fig. 3 shows the fraction of the observed photons that underwent 0, 1, 2, 3, etc. scatterings before reaching the observer for eight values of $`\theta _0`$ as discussed above (cf. Table 1), with the solid and dotted lines showing the results for Solar and twice Solar abundance of iron, respectively. The fraction of photons which arrive without being scattered is 19% and 63% respectively for a “thick” ($`\theta _0=80^{}`$) and “skinny” ($`\theta _0=10^{}`$) torus. For $`2\times `$ Solar abundance of Fe these numbers are 32% and 75%. The data in Fig. 2 imply that fewer than 40% of the observed photons are scattered over path lengths longer than 1 light day, so the half-angle subtended by the optically thick material is less than $`10^{}`$. This implies a rather small scale height, and perhaps is due to the same material which produces the H<sub>2</sub>O maser emission (Greenhill et al. 1997).
The details of the geometry of Seyfert 2s are important in the assessment of their contribution to the Cosmic X–ray background, as the heavily absorbed AGN were postulated to make up the bulk of it (cf. Krolik et al. 1994; Madau, Ghisellini, & Fabian 1994; Comastri et al. 1995). The value of $`\theta _0`$ of $`10^{}`$ is marginally consistent with the torus geometry inferred from recent observations at hard X–ray energies. These observations show that Seyfert 2s outnumber Seyfert 1s by a factor of 4:1, while more than a quarter of Seyfert 2s have a column which is optically thick (see Giommi et al. 1998, and Gilli, Risaliti, & Salvati 1999 and references therein). Assuming a single geometry for the Seyfert nuclei where the torus has a rectangular cross-section, we can attempt to reproduce these ratios by assuming that if the central flux barely “grazes” the torus, we classify the object as any Seyfert 2, while an object is an optically thick Seyfert 2 only if the line of sight encounters the entire radial distance in the torus. In this context, requiring such 1:4:1 ratio would then imply that the torus is somewhat flattened, with outer radius of 6.5 $`\times `$ its inner radius and equatorial height of 1.3 $`\times `$ its inner radius. In this scenario, all Seyfert 2s are then seen at angles smaller than $`50^{}`$ from the plane, while the optically thick Seyfert 2s are confined to angles of $`12^{}`$. We repeated the Monte–Carlo calculations with this rectangular geometry and find that the fraction of scattered photons is $``$ 50%, still too large to match the observed hard X–ray variability. A further problem arises if this is indeed a universal geometry for all Seyferts as the Thomson depth of the absorber is large, $`1.7`$ in the more restrictive $`2\times `$ Solar case. This absorber reduces the unabsorbed flux $``$ five-fold or more, and this is not consistent with the 1:4:1 ratio, by a large margin. We thus conclude that a population of AGN with geometry very nearly that of NGC 4945 cannot make up the CXB. Instead, significantly larger fraction of the heavily obscured AGN is required (implying a large solid angle subtended by the absorber) than implied from the rather small $`\theta _0`$ inferred by us. One plausible scenario would have the local optically thick Seyfert 2s surrounded by absorbers that already collapsed to an accretion disk, while in the more distant objects – in the earlier stages of evolution – such absorbers had larger vertical extent.
Alternately, the absorption could come from a structure which is much closer to the nucleus. The variability limits impose constraints on the amount of scattered flux that is lagged on time scales of more than 1 day, and thus they do not constrain the height of any structure which is $`<<1`$ light day from the X–ray source (corresponding to a distance of $`<<10^4`$ Schwarzschild radii for the mass of the black hole of $`10^6`$ $`M_{}`$). While a very geometrically thick accretion disk cannot be ruled out from our data, there are considerable theoretical difficulties in maintaining a structure with a large height scale. It is far easier to envisage a structure with a small height scale such an accretion disk with outer regions somewhat thickened by instabilities resulting from radiation pressure warping (cf. Maloney, Begelman, & Pringle 1996).
With these arguments for the Thomson-thick absorber subtending a small solid angle in NGC 4945, we can now estimate the true luminosity of the source. The Monte Carlo simulations show that the intrinsic flux must be substantially larger than that observed (by a factor of $`e^{\tau _\mathrm{e}}`$, or $`\times 11`$ for Solar abundances) because any scattered nuclear photons are lost from the line of sight. The 1-500 keV flux, corrected for photoelectric absorption alone, is $`5\times 10^{10}`$ erg cm<sup>-2</sup> s<sup>-1</sup>, so correcting for the Thomson opacity yields a 1-500 keV intrinsic X–ray luminosity of $`10^{43}`$ erg s<sup>-1</sup>. Assuming that the thermal (opt/UV/EUV) emission from the accretion disk is roughly equal to the hard X–ray emission gives a total bolometric luminosity of the nucleus of $`2\times 10^{43}`$ erg s<sup>-1</sup>. With this, and $`M_{\mathrm{BH}}`$ of $`1.4\times 10^6`$ $`M_{}`$, the source is radiating at $``$ 10% of the Eddington luminosity; even if the abundances are twice-Solar (which yields $`\tau _\mathrm{e}`$ of 1.7), $`L/L_{\mathrm{Edd}}`$ is at least $``$ 5%. (We note that similar $`L/L_{\mathrm{Edd}}`$ was also inferred by Greenhill et al. (1997), but the discovery of the rapid hard X–ray variability allows us to determine the source luminosity more accurately, as now we know that relatively few photons are scattered back to the line of sight.) NGC 4945 is one of the few AGN where this quantity can be calculated robustly, since the mass of the central object is known, although we are aware that the value of its central mass is not as accurate as for the famous megamaser NGC 4258; the resulting uncertainty in the estimate of $`L/L_{\mathrm{Edd}}`$ may be a factor of 2, comparable to the effects of the unknown Fe abundance or the ratio of $`L_{\mathrm{Tot}}/L_{\mathrm{X}\mathrm{ray}}`$. The resulting $`L/L_{\mathrm{Edd}}`$ is comparable to that inferred for the well studied Seyfert 2 NGC 1068, although since the absorber is completely opaque even to hard X–rays, the central luminosity can be estimated only indirectly. These two AGN are at the opposite end of the scale to NGC 4258, which radiates at $`10^4`$ $`L_{\mathrm{Edd}}`$ or less (Lasota et al. 1996). The mass accretion rates inferred from those values of $`L/L_{\mathrm{Edd}}`$ put strong constraints on possible underlying radiation mechanisms. While the recently popular advective disk models can produce X–ray hot flows up to about 10% of the Eddington limit, these collapse at higher $`L/L_{\mathrm{Edd}}`$: our data show that NGC 4945 lies perilously close to this limits.
ACKNOWLEDGEMENTS: We thank the RXTE satellite team for scheduling the observations allowing the daily sampling, Tess Jaffe for her help with the RXTE data reduction via her indispensable script rex, and Dr. Julian Krolik for his helpful comments on the manuscript. This project was partially supported by ITP/NSF grant PHY94-07194, NASA grants and contracts to University of Maryland and USRA, and the Polish KBN grant 2P03D01816.
Figure Captions
Fig. 1: Broad-band unfolded X–ray spectrum of NGC 4945 as measured with the RXTE PCA and HEXTE instruments. The data were fitted with a phenomenological model which includes a hard power law component photo–electrically absorbed by neutral gas with Solar abundances at a column of $`4.5\pm 0.3\times 10^{24}`$ cm<sup>-2</sup>, with photon spectral index $`\mathrm{\Gamma }=1.45_{0.1}^{+0.1}`$, exponentially cutting off at 100 keV, plus a non-variable soft component (assumed to be a power law), and a Fe K line. The observed 8 – 30 keV flux of the hard component is $`5\times 10^{11}`$ erg cm<sup>-2</sup> s<sup>-1</sup>.
Fig. 2: Hard X–ray light curve of the Seyfert 2 galaxy NGC 4945 measured with the RXTE PCA instrument, showing a rapid, large amplitude flux variability. Plotted are data from all three layers of three PCA detectors that were turned on during all pointings, over the energy channels nominally corresponding to the range of 8 – 30 keV.
Fig. 3: Fraction of the observed photons reaching an observer located in the equatorial plane of a torus plotted against the number of scatterings that those photons encountered before reaching an observer. The angle given in each panel is the vertical half-angle $`\theta _0`$ subtended by the torus as seen from the central source. Iron abundance (relative to Solar) is assumed to be 1 (solid line) and 2 (dotted line), with $`\tau _\mathrm{e}`$ equal to the best fit value for a given Fe abundance and $`\theta _0`$ (cf. Table 1). Since the fractional amplitude of variability on short time scales (cf. Fig. 2) is large ($`>`$ 60%), $`\theta _0`$ of the optically thick structure must be small, so that majority of the photons reaching an observer are not scattered. |
warning/0002/astro-ph0002203.html | ar5iv | text | # Nuclear Matter and its Role in Supernovae, Neutron Stars and Compact Object Binary Mergers
## 1 Introduction
The equation of state (EOS) of dense matter plays an important role in the supernova phenomenon and in the structure and evolution of neutron stars. Matter in the collapsing core of a massive star at the end of its life is compressed from white dwarf-like densities of about $`10^6`$ g cm<sup>-3</sup> to two or three times the nuclear saturation density, about $`310^{14}`$ g cm<sup>-3</sup> or $`n_s=0.16`$ baryons fm<sup>-3</sup>. The central densities of neutron stars may range up to 5–10 $`n_s`$. At densities around $`n_s`$ and below matter may be regarded as a mixture of neutrons, protons, electrons and positrons, neutrinos and antineutrinos, and photons. At higher densities, additional constituents, such as hyperons, kaons, pions and quarks may be present, and there is no general consensus regarding the properties of such ultradense matter. Fortunately for astrophysics, however, the supernova phenomenon and many aspects of neutron star structure may not depend upon ultradense matter, and this article will focus on the properties of matter at lower densities.
The main problem is to establish the state of the nucleons, which may be either bound in nuclei or be essentially free in continuum states. Neither temperatures nor densities are large enough to excite degrees of freedom, such as hyperons, mesons or quarks. Electrons are rather weakly interacting and may be treated as an ideal Fermi gas: at densities above $`10^7`$ g cm<sup>-3</sup>, they are relativisitic. Because of their even weaker interactions, photons and neutrinos (when they are confined in matter) may also be treated as ideal gases.
At low enough densities and temperatures, and provided the matter does not have too large a neutron excess, the relevant nuclei are stable in the laboratory, and experimental information may be used directly. The so-called Saha equation may be used to determine their relative abundances. Under more extreme conditions, there are a number of important physical effects which must be taken into account. At higher densities, or at moderate temperatures, the neutron chemical potential increases to the extent that the density of nucleons outside nuclei can become large. It is then important to treat matter outside nuclei in a consistent fashion with that inside. These nucleons will modify the nuclear surface, decreasing the surface tension. At finite temperatures, nuclear excited states become populated, and these states can be included by treating nuclei as warm drops of nuclear matter. At low temperatures, nucleons in nuclei are degenerate and Fermi-liquid theory is probably adequate for their description. However, near the critical temperature above which the dense phase of matter inside nuclei can no longer coexist with the lighter phase of matter outside nuclei, the equilibrium of the two phases of matter is crucial.
The fact that at subnuclear densities the spacing between nuclei may be of the same order of magnitude as the nuclear size itself will lead to substantial reductions in the nuclear Coulomb energy. Although finite temperature “plasma” effects will modify this, the zero-temperature Wigner-Seitz approximation employed by Baym, Bethe & Pethick BBP is usually adequate. Near the nuclear saturation density, nuclear deformations must be dealt with, including the possibilities of “pasta-like” phases and matter turning “inside-out” (i.e., the dense nuclear matter envelopes a lighter, more neutron-rich, liquid). Finally, the translational energy of the nuclei may be important under some conditions. This energy is important in that it may substantially reduce the average size of the nuclear clusters.
An acceptable way of bridging the regions of low density and temperature, in which the nuclei can be described in terms of a simple mass formula, and high densities and/or high temperatures in which the matter is a uniform bulk fluid, is to use a compressible liquid droplet model for nuclei in which the drop maintains thermal, mechanical, and chemical equilibrium with its surroundings. This allows us to address both the phase equilibrium of nuclear matter, which ultimately determines the densities and temperatures in which nuclei are permitted, and the effects of an external nucleon fluid on the properties of nuclei. Such a model was originally developed by Lattimer et al. LPRL and modified by Lattimer & Swesty LS . This work was a direct result of David Schramm’s legendary ability to mesh research activities of various groups, in this case to pursue the problem of neutron star decompression. After the fact, the importance of this topic for supernovae became apparent.
## 2 Nucleon Matter Properties
The compressible liquid droplet model rests upon the important fact that in a many-body system the nucleon-nucleon interaction exhibits saturation. Empirically, the energy per particle of bulk nuclear matter reaches a minimum, about –16 MeV, at a density $`n_s0.16\mathrm{fm}^3`$. Thus, close to $`n_s`$, its density dependence is approximately parabolic. The nucleon-nucleon interaction is optimized for equal numbers of neutrons and protons (symmetric matter), so a parabolic dependence on the neutron excess or proton fraction, $`x`$, can be assumed. About a third to a half of the energy change made by going to asymmetric matter is due to the nucleon kinetic energies, and to a good approximation, this varies as $`(12x)^2`$ all the way to pure neutron matter ($`x=0`$). The $`x`$ dependence of the potential terms in most theoretical models can also be well approximated by a quadratic dependence. Finally, since at low temperatures the nucleons remain degenerate, their temperature dependence to leading order is also quadratic. Therefore, for analytical purposes, the nucleon free energy per baryon can be approximated as $`f_{bulk}(n,x)`$, in MeV, as
$`f_{bulk}(n,x)16+S_v(n)(12x)^2+{\displaystyle \frac{K_s}{18}}\left({\displaystyle \frac{n}{n_s}}1\right)^2`$ (1)
$`{\displaystyle \frac{K_s^{}}{27}}\left({\displaystyle \frac{n}{n_s}}1\right)^3a_v(n,x)T^2,`$ (2)
where $`a_v(n)=(2m^{}/\mathrm{}^2)(\pi /12n)^{2/3}`$. The expansion parameters, whose values are uncertain to varying degrees, are the incompressibility, $`K_s=190250`$ MeV, the skewness parameter $`K_s^{}=17802380`$ MeV, the symmetry energy coefficient $`S_vS_v(n_s)=2536`$ MeV, and the bulk level density parameter, $`a_v(n_s,x=1/2)(1/15)(m^{}(n_s,x=1/2)/m)`$ MeV<sup>-1</sup>, where $`m^{}`$ is the effective mass of the nucleon. Values for $`m^{}(n_s,x=1/2)/m`$ are in the range $`0.70.9`$. The general definition of the incompressibility is $`K=9dP/dn=9d(n^2df_{bulk}/dn)/dn`$, where $`P`$ is the pressure, and $`K_sK(n_s,1/2)`$. It is worthwhile noting that the symmetry energy and nucleon effective mass (which directly affects the matter’s specific heat) are density dependent, but these dependencies are difficult to determine from experiments. The parameters, and their density dependences, characterize the nuclear force model and are essential to our understanding of astrophysical phenomena.
The experimental determination of these parameters has come from comparison of the total masses and energies of giant resonances of laboratory nuclei with theoretical predictions. Some of these comparisons are easily illustrated with the compressible liquid droplet model. In this model, the nucleus is treated as uniform drop of nuclear matter with temperature $`T`$, density $`n_i`$ and proton fraction $`x_i`$. The nucleus will, in general, be surrounded by and be in equilibrium with a vapor of matter with density $`n_o`$ and proton fraction $`x_o`$. At low ambient densities $`n`$ and vanishing temperature, the outside vapor vanishes. Even at zero temperature, if $`n`$ is large enough, greater than the so-called neutron drip density $`n_d1.610^3`$ fm<sup>-3</sup>, the neutron chemical potential of the nucleus is positive and “free” neutrons exist outside the nucleus. At finite temperature, the external vapor consists of both neutrons and protons. In addition, because of their high binding energy, $`\alpha `$particles will also be present. The total free energy density is the sum of the various components:
$$F=F_H+F_o+F_\alpha +F_e+F_\gamma .$$
(3)
Here, $`F_H`$ and $`F_o`$ represent the free energy densities of the heavy nuclei and the outside vapor, respectively. The energy densities of the electrons and photons, $`F_e`$ and $`F_\gamma `$, are independent of the baryons and play no role in the equilibrium. For simplicity, we neglect the role of $`\alpha `$-particles in the following discussion (although it is straightforward to include their effect LPRL ).
In the compressible liquid drop model, it is assumed that the nuclear energy can be written as an expansion in $`A^{1/3}`$ and $`(12x_i)^2`$:
$$F_H=un_i[f_{bulk}+f_{surf}+f_{Coul}+f_{trans}],$$
(4)
where the $`f`$’s represent free energies per baryon due to the bulk, surface, Coulomb, and translation, respectively. The bulk energy, for example, is given by Eq. (2). The surface energy can be parametrized as
$$f_{surf}=4\pi R^2\sigma (x_i,T)4\pi R^2h(T)[\sigma _o\sigma _s(12x_i)^2],$$
(5)
where $`R`$ is the nuclear charge radius, $`h(T)`$ is a calculable function of temperature, $`\sigma _o`$ is the surface tension of symmetric matter, and $`\sigma _s=(n_i^2/36\pi )^{1/3}S_s`$ where $`S_s`$ is the surface symmetry energy coefficient from the traditional mass formula. In this simplified discussion, the influence of the neutron skin LPRL , which distinguishes the “drop model” from the “droplet model”, is omitted. The Coulomb energy, in the Wigner-Seitz approximation BBP , is
$$f_{Coul}=0.6x_i^2A^2e^2D(u)/R,$$
(6)
where $`D(u)=11.5u^{1/3}+0.5u`$ and $`u`$ is the fraction of the volume occupied by nuclei. If the fractional mass of matter outside the nuclei is small, $`un/n_i`$.
It is clear that additional parameters, $`S_s`$ and another involving the temperature dependence of $`h`$, exist in conjunction with those defining the expansions of the bulk energy. The temperature dependence is related to the matter’s critical temperature $`T_c`$ at which the surface disappears. It is straightforward to demonstrate from the thermodynamic relations defining $`T_c`$, namely $`P_{bulk}/n=0`$ and $`^2P_{bulk}/n^2=0`$, that $`T_c\sqrt{K_s}`$. Therefore, the specific heat to be associated with the surface energy will in general be proportional to $`T_c^2K_s^1`$. About half the total specific heat originates in the surface, so $`K_s`$ influences the temperature for a given matter entropy, important during stellar collapse.
The equilibrium between nuclei and their surroundings is determined by minimizing $`F`$ with respect to its internal variables, at fixed $`n,Y_e`$, and $`T`$. This is described in more detail in Refs. LPRL ; LS , and leads to equilibrium conditions involving the pressure and the baryon chemical potentials, as well as a condition determining the nuclear size $`R`$. The latter is analogous to the one found by Baym, Bethe & Pethick BBP who equated the nuclear surface energy with twice the Coulomb energy. The relations in Eqs. (5) and (6) lead to
$$R=\left[\frac{15\sigma (x_i)}{8\pi e^2x_i^2n_i^2}\right]^{1/3}.$$
(7)
Experimental limits to $`K_s`$, most importantly from RPA analyses of the breathing mode of the giant monopole resonance Blaizot , give $`K_s230`$ MeV. It is also possible to obtain values from the so-called scaling model developed from the compressibile liquid drop model. The finite-nucleus incompressibility is
$$K(A,Z)=(M/\mathrm{}^2)R^2E_{br}^2,$$
(8)
where $`M`$ is the mass of the nucleus and $`E_{br}`$ is the breathing-mode energy. $`K(A,Z)`$ is commonly expanded as
$$K(A,Z)=K_s+K_{surf}A^{1/3}+K_{vI}I^2+K_{surfI}I^2A^{1/3}+K_CZ^2A^{4/3},$$
(9)
and then fit by least squares to the data for $`E_{br}`$. Here the asymmetry $`I=12Z/A`$. For a given assumed value of $`K_s`$, and taking $`K_{surfI}=0`$, Pearson P showed that experimental data gave
$`K_C15.40.065K_s\pm 2\mathrm{MeV},K_{surf}2303.2K_s\pm 50\mathrm{MeV}.`$ (10)
With minimal assumptions regarding the form of the nuclear force, Pearson P demonstrated that values of $`K_s`$ ranging from 200 MeV to more than 350 MeV could be consistent with experimental data.
But the liquid drop model predicts other relations between the parameters:
$`K(A,Z)`$ $`=`$ $`R^2{\displaystyle \frac{^2E(Z,A)/A}{R^2}}|_A=9n^2{\displaystyle \frac{^2E(Z,A)/A}{n^2}}|_A,`$ (11)
$`0=P(A,Z)`$ $`=`$ $`R{\displaystyle \frac{^2E(Z,A)/A}{R}}|_A=3n{\displaystyle \frac{E(Z,A)/A}{n}}|_A.`$ (12)
Here $`E(Z,A)`$ is the total energy of the nucleus, and is equivalent to Eq. (4). The second of these equations simply expresses the equilibrium between the nucleus and the surrounding vacuum, which implies that the pressure of the bulk matter inside the nucleus is balanced by the pressure due to the curvature of the surface and the Coulomb energy. It can then be shown that
$`K_C`$ $`=`$ $`(3e^2/5r_o)[8+27n_s^3f_{bulk}^{\prime \prime \prime }(n_s)/K_s],`$ (13)
$`K_{surf}`$ $`=`$ $`4\pi r_o^2\sigma _o[9n_s^2\sigma _o^{\prime \prime }/\sigma _o+22+54n_s^3f_{bulk}^{\prime \prime \prime }(n_s)/K_s],`$ (14)
$`K_{surfI}`$ $`=`$ $`4\pi r_o^2\sigma _s[9n_s^2\sigma _s^{\prime \prime }/\sigma _s+22+54n_s^3f_{bulk}^{\prime \prime \prime }(n_s)/K_s],`$ (15)
$`K_I`$ $`=`$ $`9[n_s^2S_v^{\prime \prime }(n_s)2n_oS_v^{}(n_s)9n_s^4S_v^{}(n_s)f_{bulk}^{\prime \prime \prime }(n_s)/K_s].`$ (16)
Primes denote derivatives with respect to the density. From these relations, and again assuming $`K_{surfI}=0`$, Pearson demonstrated that an interesting correlation between $`K_s`$ and $`K_s^{}`$, where $`K_s^{}27n_s^3f_{bulk}^{\prime \prime \prime }(n_s)`$, could be obtained:
$$K_s^{}=0.0860K_s^2+(28.37\pm 2.65)K_s.$$
(17)
Assuming $`K_s190250`$ MeV, this suggests that $`K_s^{}=17802380`$ MeV, a potential constraint. Alternatively, eliminating $`K_s^{}`$, one finds
$$K_s=137.426.36n_s^2\sigma _o^{\prime \prime }/\sigma _o\pm 23.2\mathrm{MeV}.$$
(18)
The second derivative of the surface tension can be deduced from Hartree-Fock or Thomas-Fermi semi-infinite surface calculations. For example, if a parabolic form of $`f_{bulk}`$ is used, one finds
$$n_s^2\sigma _o^{\prime \prime }/\sigma _o=6$$
(19)
leading to $`K_s=295.5\pm 23.2`$ MeV. In general, the density dependence of $`S_v`$ will decrease the magnitudes of $`K_s`$ and $`\sigma _o^{\prime \prime }`$ from the above values.
It is hoped current experimental work will tighten these constraints. A shortcoming of the scaling model is that, to date, the surface symmetry energy term was neglected. This is not required, however, and further work is necessary to resolve this matter.
Because the surface energy represents the energy difference between uniformly and realistically distributed nuclear material in a nucleus, the parameter $`S_s`$ can be related to the density dependence of $`S_v(n)`$ and to $`K_s`$. If $`f_{bulk}`$ is assumed to behave quadratically with density around $`n_s`$, this relation can be particularly simply expressed L :
$$\frac{S_s}{S_v}=\frac{3}{\sqrt{2}}\frac{a_{1/2}}{r_o}_0^1\frac{\sqrt{x}}{1x}\left[\frac{S_v}{S_v(xn_s)}1\right]𝑑x.$$
(20)
Here, $`S_vS_v(n_s)`$, $`a_{1/2}=(dr/d\mathrm{ln}n)_{n_s/2}`$ is a measure of the thickness of the nuclear surface and $`r_o=(4\pi n_s/3)^{1/3}=R/A^{1/3}`$. If $`S_v(n)`$ is linear, then the integral is 2; if $`S_v(n)n^{2/3}`$, then the integral is 0.927. Since $`a_{1/2}`$ will be sensitive to the value of $`K_s`$, we expect the value of $`S_s/S_v`$ to be also.
Experimentally, there are two major sources of information regarding the symmetry energy parameters: nuclear masses and giant resonance energies. However, because of the small excursions in $`A^{1/3}`$ afforded by laboratory nuclei, each source provides only a correlation between $`S_s`$ and $`S_v`$. For example, the total symmetry energy in the liquid droplet model (now explicitly including the presence of the neutron skin, see Ref. LPRL ) is
$$E_{sym}=(12x_i)^2S_v/[1+(S_s/S_v)A^{1/3}].$$
(21)
Evaluating $`\alpha =d\mathrm{ln}S_s/d\mathrm{ln}S_v`$ near the “best-fit” values $`S_{s0}`$ and $`S_{v0}`$, one finds
$$\alpha 2+S_{v0}<A>^{1/3}/S_{s0}6,$$
(22)
where $`<A>^{1/3}`$ for the fitted nuclei is about 5. Thus, as the value of $`S_v`$ is changed in the mass formula, the value of $`S_s`$ must vary rapidly to compensate.
An additional correlation between these parameters can be obtained from the fitting of isovector giant resonances, and this has the potential of breaking the degeneracy of $`S_v`$ and $`S_s`$, because it has a different slope L . Lipparini & Stringari Lip used a hydrodynamical model of the nucleus to derive the isovector resonance energy:
$`E_d`$ $`=`$ $`\sqrt{{\displaystyle \frac{24\mathrm{}^2}{m^{}}}{\displaystyle \frac{NZ}{A}}\left[{\displaystyle \frac{nr^2S_v}{S_v(n)}d^3r}\right]^1}`$ (23)
$``$ $`96.5\sqrt{{\displaystyle \frac{m}{m^{}}}{\displaystyle \frac{S_v}{30\mathrm{MeV}}}\left[1+{\displaystyle \frac{5S_s}{3S_vA^{1/3}}}\right]^1}A^{1/3}\mathrm{MeV},`$ (24)
where $`m^{}`$ is an effective nucleon mass. This relation results in a slightly less-steep correlation between $`S_s`$ and $`S_v`$,
$$\alpha =2/m^{}+(3/5)S_{v0}<A>^{1/3}/S_{s0}45.$$
(25)
Unfortunately the value of $`m^{}`$ is an undetermined parameter and this slope is not very different from that obtained from fitting masses. Therefore, uncertainties in the model make a large difference to the crossing point of these two correlations. A strong theoretical attack, perhaps using further RPA analysis, together with more experiments to supplement the relatively meager amount of existing data, would be very useful.
## 3 The Equation of State and the Collapse of Massive Stars
Massive stars at the end of their lives are believed to consist of a white dwarf-like iron core of 1.2–1.6 M having low entropy ($`s1`$), surrounded by layers of less processed material from shell nuclear burning. The effective Chandrasekhar mass, the maximum mass the degenerate electron gas can support, is dictated by the entropy and the average lepton content, $`Y_L`$, believed to be around 0.41–0.43. As mass is added to the core by shell Si-burning, the core eventually becomes unstable and collapses.
During the collapse, the lepton content decreases due to net electron capture on nuclei and free protons. But when the core density approaches $`10^{12}`$ g cm<sup>-3</sup>, the neutrinos can no longer escape from the core on the dynamical collapse time Sato . After neutrinos become trapped, $`Y_L`$ is frozen at a value of about 0.38–0.40, and the entropy is also thereafter fixed. The core continues to collapse until the rapidly increasing pressure reverses the collapse at a bounce density of a few times nuclear density.
The immediate outcome of the shock generated by the bounce is also dependent upon $`Y_L`$. First, the shock energy is determined by the net binding energy of the post-bounce core, and is proportional to $`Y_e^{10/3}`$ LBY . Second, the shock is largely dissipated by the energy required to dissociate massive nuclei in the still-infalling matter of the original iron core outside the post-bounce core. The larger the $`Y_L`$ of the core, the larger its mass and the smaller this shell. Therefore, the progress of the shock is very sensitive to the value of $`Y_L`$.
The final value of $`Y_L`$ is controlled by weak interaction rates, and is strongly dependent upon the fraction of free protons, $`X_p`$, which is proportional to $`\mathrm{exp}(\mu _p/T)`$, and the phase space available for proton capture on nuclei, which is proportional to $`\mu _e\widehat{\mu }`$, where $`\widehat{\mu }=\mu _n\mu _p`$. Both are sensitive to the proton fraction in nuclei ($`x_i`$) and are largely controlled by $`Y_L`$. In addition, the specific heat controls the temperature which has a direct influence upon the free proton abundance and the net electron capture rate. In spite of the intricate feedback, nuclear parameters relating chemical potentials to composition, especially $`S_v`$ and $`S_s`$, are obviously important.
As an example, consider $`\widehat{\mu }=\mu _n\mu _p=n_{i}^{}{}_{}{}^{1}F_H/x_i`$. With the model of Eqs. (4)-(6), one has
$$\widehat{\mu }=4S_v(12x_i)\left(\frac{72\pi e^2D}{5x_in_i}\right)^{1/3}\frac{\sigma _o\sigma _s(12x_i)(16x_i)}{(\sigma _o\sigma _s(12x_i)^2)^{1/3}}.$$
(26)
Recall that $`\sigma _sS_s`$. Although the bulk and Coulomb terms alone (Eq. 26 with $`\sigma _s=0`$) imply that $`\widehat{\mu }`$ for a given $`x_i`$ rises with increasing $`S_v`$, the proper inclusion of the surface symmetry energy gives rise to the opposite behavior. This is illustrated in Fig. 1.
Uncertainties in nuclear parameters can thus be expected to have an influence upon the collapse of massive stars, for example, in the collapse rate, the final trapped lepton fraction, and the radius at which the bounce-generated shock initially stalls. Swesty, Lattimer & Myra SLM investigated the effects upon stellar collapse of altering parameters in a fashion constrained by nuclear systematics. They found that as long as the parameters permitted a neutron star maximum mass above the PSR1913+16 mass limit (1.44 M), the shock generated by core bounce consistently stalls near 100 km, independently of the assumed $`K_s`$ in the range 180–375 MeV and $`S_v`$ in the range 27–35 MeV. Ref. SLM also found that the final trapped lepton fraction is also apparently independent of variations in both $`K_s`$ and $`S_v`$. These results are in contrast to earlier simulations which had used EOSs that could not support cold, catalyzed 1.4 M stars, or in which $`S_s`$ was not varied consistently with $`S_v`$. The strong feedback between the EOS, weak interactions, neutrino transport, and hydrodynamics is an example of Mazurek’s Law.
In fact, the only significant consequence of varying $`S_v`$ involved the pre-bounce neutrino luminosities. Increasing $`S_v`$ increases the electron capture rate (proportional to $`\mu _e\widehat{\mu }`$ and therefore increases the $`\nu _e`$ luminosity during collapse, as shown in Fig. 2. Nevertheless, the collapse rate also increases, so that neutrino trapping occurs sooner and the final trapped lepton fraction does not change. It is possible that large neutrino detectors such as Super-Kamiokande or SNO may be able to observe an enhanced early rise in neutrino luminosity from nearby galactic supernovae.
## 4 The Structure of Neutron Stars
The theoretical study of the structure of neutron stars is crucial if new observations of masses and radii are to lead to effective constraints on the EOS of dense matter. This study becomes ever more important as laboratory studies may be on the verge of yielding evidence about the composition and stiffness of matter beyond $`n_s`$. To date, several accurate mass determinations of neutron stars are available, and they all lie in a narrow range ($`1.251.44`$ M). There is some speculation that the absence of neutron stars with masses above 1.5 M implies that $`M_{max}`$ for neutron stars has approximately this value. However, since fewer than 10 neutron stars have been weighed, and all these are in binaries, this conjecture is premature. Theoretical studies of dense matter indicate that considerable uncertainties exist in the high-density behavior of the EOS largely because of the poorly constrained many-body interactions. These uncertainties are reflected in a significant uncertainty in the maximum mass of a beta-stable neutron star, which ranges from 1.5–2.5 M.
There is some theoretical support for a lower mass limit for neutron stars in the range $`1.11.2`$ M. This follows from the facts that the collapsing core of a massive star is always greater than 1 M and the minimum mass of a protoneutron star with a low-entropy inner core of $`0.6`$ M and a high-entropy envelope is at least 1.1 M. Observations from the Earth of thermal radiation from neutron star surfaces could yield values of the quantity $`R_{\mathrm{}}=R/\sqrt{12GM/Rc^2}`$, which results from redshifting the stars luminosity and temperature. $`MR`$ trajectories for representative EOSs (discussed below) are shown in Figure 3. It appears difficult to simultaneously have $`M>1`$M and $`R_{\mathrm{}}<12`$ km. Those pulsars with at least some suspected thermal radiation generically yield effective values of $`R_{\mathrm{}}`$ so small that it is believed that the radiation originates from polar hot spots rather than from the surface as a whole. Other attempts to deduce a radius include analyses Tit of X-ray bursts from sources 4U 1705-44 and 4U 1820-30 which implied rather small values, $`9.5<R_{\mathrm{}}<14`$ km. However, the modeling of the photospheric expansion and touchdown on the neutron star surface requires a model dependent relationship between the color and effective temperatures, rendering these estimates uncertain. Absorption lines in X-ray spectra have also been investigated with a view to deducing the neutron star radius. Candidates for the matter producing the absorption lines are either the accreted matter from the companion star or the products of nuclear burning in the bursts. In the former case, the most plausible element is thought to be Fe, in which case the relation $`R3.2GM/c^2`$, only slightly larger than the minimum possible value based upon causality, LPMY ; glen is inferred. In the latter case, plausible candidates are Ti and Cr, and larger values of the radius would be obtained. In both cases, serious difficulties remain in interpreting the large line widths, of order 100–500 eV, in the $`4.1\pm 0.1`$ keV line observed from many sources. A first attempt at using light curves and pulse fractions from pulsars to explore the $`MR`$ relation suggested relatively large radii, of order 15 km Page . However, this method, which assumed dipolar magnetic fields, was unable to satisfactorily reconcile the calculated magnitudes of the pulse fractions and the shapes of the light curves with observations.
Prospects for a radius determination have improved in recent years, however, with the detection of a nearby neutron star, RX J185635-3754, in X-rays and optical radiation Walter . The observed X-rays, from the ROSAT satellite, are consistent with blackbody emission with an effective temperature of about 57 eV and very little extinction. In addition, the fortuitous location of the star in the foreground of the R CrA molecular cloud limits the distance to $`D<120`$ pc. The fact that the source is not observable in radio and its lack of variability in X-rays implies that it is not a pulsar unlike other identified radio-silent isolated neutron stars. This gives the hope that the observed radiation is not contaminated with non-thermal emission as is the case for pulsars. The X-ray observations of RXJ185635-3754 alone yield $`R_{\mathrm{}}7.3(D/120\mathrm{pc})\mathrm{km}`$ for a best-fit blackbody. Such a value is too small to be consistent with any neutron star with more than 1 M. But the optical flux is about a factor of 2.5 brighter than what is predicted for the X-ray blackbody, which is consistent with there being a heavy-element atmosphere Romani . With such an atmosphere, it is found ALPW that the effective temperature is reduced to approximately 50 eV and $`R_{\mathrm{}}`$ is also increased, to a value of approximately $`21.6(D/120\mathrm{pc})\mathrm{km}`$. Upcoming parallax measurements with the Hubble Space Telescope should permit a distance determination to about 10-15% accuracy. If X-ray spectral features are discovered with the planned Chandra and XMM space observatories, the composition of the neutron star atmosphere can be inferred, and the observed redshifts will yield independent mass and radius information. In this case, both the mass and radius of this star will be found.
Furthermore, a proper motion of 0.34 <sup>′′</sup> yr<sup>-1</sup> has been detected, in a direction that is carrying the star away from the Upper Scorpius (USco) association ALPW . With an assumed distance of about 80 pc, the positions of RX J185635-3754 and this association overlap about 800,000 years ago. The runaway OB star $`\zeta `$ Oph is also moving away from USco, appearing to have been ejected on the order of a million years ago. The superposition of these three objects is interesting, and one can speculate that this is not coincidental. If upcoming parallax measurements are consistent with a distance to RX J185635-3754 of about 80 pc, the evidence for this scenario will be strong, and a good age estimate will result.
In this section, a striking empirical relationship is noted which connects the radii of neutron stars and the pressure of matter in the vicinity of $`n_s`$. In addition, a number of analytic, exact, solutions to the general relativistic TOV equation of hydrostatic equilibrium are explored that lead to several useful approximations for neutron star structure which directly correlate observables such as masses, radii, binding energies, and moments of inertia. The binding energy, of which more than 99% is carried off in neutrinos, will be revealed from future neutrino observations of supernovae. Moments of inertia are connected with glitches observed in the spin down of pulsars, and their observations yield some interesting conclusions about the distribution of the moment of inertia within the rotating neutron star. From such comparisons, it may become easier to draw conclusions about the dense matter EOS when firm observations of neutron star radii or moments of inertia become available to accompany the several known accurate mass determinations.
### 4.1 Neutron Star Radii
The composition of a neutron star chiefly depends on the nature of strong interactions, which are not well understood in dense matter. The several possible models investigated LPMY ; physrep can be conveniently grouped into three broad categories: nonrelativistic potential models, field-theoretical models, and relativistic Dirac-Brueckner-Hartree-Fock models. In each of these approaches, the presence of additional softening components such as hyperons, Bose condensates or quark matter, can be incorporated.
Figure 3 displays the mass-radius relation for several recent EOSs (the abbreviations are explained in Table 1). Even a cursory glance indicates that in the mass range from $`11.5`$ M it is usually the case that the radius has little dependence upon the mass. The lone exception is the model GS1, in which a kaon condensate, leading to considerable softening, appears. While it is generally assumed that a stiff EOS leads to both a large maximum mass and a large radius, many counter examples exist. For example, MS3 has a relatively small maximum mass but has large radii compared to most other EOSs with larger maximum masses. Also, not all EOSs with extreme softening have small radii (viz., GS2). Nonetheless, for stars with mass greater than 1 M, only models with a large degree of softening can have $`R_{\mathrm{}}<12`$ km. Should the radius of a neutron star ever be accurately determined to satisfy $`R_{\mathrm{}}<12`$ km, a strong case can be made for the existence of extreme softening.
It is relevant that a Newtonian polytrope with $`n=1`$ has the property that the stellar radius is independent of both the mass and central density. In fact, numerical relativists have often approximated equations of state with $`n=1`$ polytropes. An $`n=1`$ polytrope has the property that the radius is proportional to the square root of the constant $`K`$ in the polytropic pressure law $`P=K\rho ^{1+1/n}`$. This suggests that there might be a quantitative relation between the radius and the pressure that does not depend upon the equation of state at the highest densities, which determines the overall softness or stiffness (and hence, the maximum mass).
To make the relation between matter properties and the nominal neutron star radius definite, Fig. 4 shows the remarkable empirical correlation which exists between the radii of 1 and 1.4 M stars and the matter’s pressure evaluated at densities of 1, 1.5 and 2 $`n_s`$. Table 1 explains the EOS symbols used in Fig. 4. Despite the relative insensitivity of radius to mass for a particular “normal” equation of state, the nominal radius $`R_M`$, which is defined as the radius at a particular mass $`M`$ in solar units, still varies widely with the EOS employed. Up to $`5`$ km differences are seen in $`R_{1.4}`$, for example, in Fig. 4. This plot is restricted to EOSs which have maximum masses larger than about 1.55 M and to those which do not have strong phase transitions (such as those due to a Bose condensate or quark matter). Such EOSs violate these correlations, especially for the case of 1.4 M stars. We emphasize that this correlation is valid only for cold, catalyzed neutron stars, i.e., it will not be valid for protoneutron stars which have finite entropies and might contain trapped neutrinos. The correlation has the form
$$R\mathrm{constant}[P(n)]^{0.230.26},$$
(27)
where $`P`$ is the total pressure inclusive of leptonic contributions evaluated at the density $`n`$. An exponent of 1/4 was chosen for display in Fig. 4, but the correlation holds for a small range of exponents about this value. The correlation is marginally tighter for the baryon density $`n=1.5n_s`$ and $`2n_s`$ cases. Thus, instead of the power 1/2 that the Newtonian polytrope relations would predict, a power of approximately 1/4 is suggested when the effects of relativity are included. The value of the constant in Eq. (27) depends upon the chosen density, and can be obtained from Fig. 4.
The exponent of 1/4 can be quantitatively understood by using a relativistic generalization of the $`n=1`$ polytrope, due to Buchdahl Buchdahl . For the EOS
$$\rho =12\sqrt{p_{}P}5P,$$
(28)
where $`p_{}`$ is a constant, there is an analytic solution to Einstein’s equations:
$`e^\nu `$ $``$ $`g_{tt}=(12\beta )(1\beta u)(1\beta +u)^1;`$ (29)
$`e^\lambda `$ $``$ $`g_{rr}=(12\beta )(1\beta +u)(1\beta u)^1(1\beta +\beta \mathrm{cos}Ar^{})^2;`$ (30)
$`8\pi PG/c^4`$ $`=`$ $`A^2u^2(12\beta )(1\beta +u)^2;`$ (31)
$`8\pi \rho G/c^2`$ $`=`$ $`2A^2u(12\beta )(1\beta 3u/2)(1\beta +u)^2;`$ (32)
$`u`$ $`=`$ $`\beta (Ar^{})^1\mathrm{sin}Ar^{};r=r^{}(1\beta +u)(12\beta )^1;`$ (33)
$`A^2`$ $`=`$ $`288\pi p_{}Gc^4(12\beta )^1;R=\pi (1\beta )(12\beta )^1A^1.`$ (34)
The free parameters of this solution are $`\beta GM/Rc^2`$ and the scale $`p_{}`$. Note that $`Rp_{}^{1/2}(1+\beta ^2/2+\mathrm{})`$, so for a given value of $`p_{}`$, the radius increases only very slowly with mass, exactly as expected from an $`n=1`$ Newtonian polytrope. It is instructive to analyze the response of $`R`$ to a change of pressure at some fiducial density $`\rho `$, for a fixed mass $`M`$. One finds
$$\frac{d\mathrm{ln}R}{d\mathrm{ln}P}|_{\rho ,M}=\frac{\frac{d\mathrm{ln}R}{d\mathrm{ln}p_{}}|_\beta \frac{d\mathrm{ln}p_{}}{d\mathrm{ln}P}|_\rho }{1+\frac{d\mathrm{ln}R}{d\mathrm{ln}\beta }|_p_{}}=\left(1\frac{5}{6}\sqrt{\frac{P}{p_{}}}\right)\frac{(1\beta )(12\beta )}{2(13\beta +3\beta ^2)}.$$
(35)
In the limit $`\beta 0,P0`$ and $`d\mathrm{ln}R/d\mathrm{ln}P1/2`$, the value characteristic of an $`n=1`$ Newtonian polytrope. Finite values of $`\beta `$ and $`P`$ render the exponent smaller than 1/2. If the stellar radius is about 15 km, $`p_{}=\pi /(288R^2)4.8510^5`$ km<sup>-2</sup>. If the fiducial density is $`\rho 1.5m_bn_s2.0210^4`$ km<sup>-2</sup> (with $`m_b`$ the baryon mass), Eq. (28) implies that $`P8.510^6`$ km<sup>-2</sup>. For $`M=1.4`$ M, the value of $`\beta `$ is 0.14, and $`d\mathrm{ln}R/d\mathrm{ln}P0.31`$. This result is mildly sensitive to the choices for $`\rho `$ and $`R`$, and the Buchdahl solution is not a perfect representation of realistic EOSs; nevertheless, it provides a reasonable explanation of the correlation in Eq. (27).
The existence of this correlation is significant because, in large part, the pressure of degenerate matter near the nuclear saturation density $`n_s`$ is determined by the symmetry properties of the EOS. Thus, the measurement of a neutron star radius, if not so small as to indicate extreme softening, could provide an important clue to the symmetry properties of matter. In either case, valuable information is obtained.
The specific energy of nuclear matter near the saturation density may be expressed as an expansion in the asymmetry $`(12x)`$, as displayed in Eq. (2), that can be terminated after the quadratic term PAL . Leptonic contributions must be added to Eq. (2) to obtain the total energy and pressure; the electron energy per baryon is $`f_e=(3/4)\mathrm{}cx(3\pi ^2nx)^{1/3}`$. Matter in neutron stars is in beta equilibrium, i.e., $`\mu _e\mu _n+\mu _p=(f_{bulk}+f_e)/x=0`$, so the electronic contributions may be eliminated to recast the pressure as Ppuri
$`P=n^2[S_v^{}(n)(12x)^2+{\displaystyle \frac{xS_v(n)}{n}}(12x)+`$ (36)
$`{\displaystyle \frac{K_s}{9n_s}}({\displaystyle \frac{n}{n_s}}1){\displaystyle \frac{K_s^{}}{54n_s}}({\displaystyle \frac{n}{n_s}}1)^2],`$ (37)
where $`x`$ is now the beta equilibrium value. At the saturation density,
$`P_s=n_s(12x_s)[n_sS_v^{}(n_s)(12x_s)+S_vx_s],`$ (38)
where the equilibrium proton fraction at $`n_s`$ is
$`x_s(3\pi ^2n_s)^1(4S_v/\mathrm{}c)^30.04`$ (39)
for $`S_v=30`$ MeV. Due to the small value of $`x_s`$, one finds that $`P_sn_s^2S_v^{}(n_s)`$. If the pressure is evaluated at a larger density, other nuclear parameters besides $`S_v`$ and $`S_v^{}(n_s)`$, become significant. For $`n=2n_s`$, one thus has
$`P(2n_s)4n_s[n_sS_v^{}(2n_s)+(K_sK_s^{}/6)/9].`$ (40)
If it is assumed that $`S_v(n)`$ is linear in density, $`K_s220`$ MeV and $`K_s^{}2000`$ MeV (as indicated in Eq. 17), the symmetry contribution is still about 70% of the total.
The sensitivity of the radius to the symmetry energy is graphically shown by the parametrized EOS of PAL PAL in Fig. 5. The symmetry energy function $`S_v(n)`$ is a direct input in this parametrization. The figure shows the dependence of mass-radius trajectories as the quantities $`S_v`$ and $`S_v(n)`$ are alternately varied. Clearly, the density dependence of $`S_v(n)`$ is more important in determining the neutron star radius. Note also the weak sensitivity of the maximum neutron star mass to $`S_v`$.
At present, experimental guidance concerning the density dependence of the symmetry energy is limited and mostly based upon the division of the nuclear symmetry energy between volume and surface contributions, as discussed in the previous section. Upcoming experiments involving heavy-ion collisions (at GSI, Darmstadt), which might sample densities up to $`(34)n_s`$, will be limited to analyzing properties of the symmetric nuclear matter EOS through a study of matter, momentum, and energy flow of nucleons. Thus, studies of heavy nuclei far off the neutron drip lines will be necessary in order to pin down the properties of the neutron-rich regimes encountered in neutron stars.
### 4.2 Neutron Star Moments of Inertia and Binding Energies
Besides the stellar radius, other global attributes of neutron stars are potentially observable, including the moment of inertia and the binding energy. These quantities depend primarily upon the ratio $`M/R`$ as opposed to details of the EOS, as can be readily seen by evaluating them using analytic solutions to Einstein’s equations. Although over 100 analytic solutions to Einstein’s equations are known Delgaty , nearly all of them are physically unrealistic. However, three analytic solutions are of particular interest in neutron star structure.
The first is the well-known Schwarzschild interior solution for an incompressible fluid, $`\rho =\rho _c`$, where $`\rho `$ is the mass-energy density. This is mostly of interest because it determines the maximum compression $`\beta =GM/Rc^2`$ for a neutron star, namely 4/9, based upon the pressure being finite. Two aspects of the incompressible fluid that are physically unrealistic, however, include the fact that the sound speed is everywhere infinite, and that the density does not vanish on the star’s surface.
The second analytic solution, B1, due to Buchdahl Buchdahl , is described in Eq. (34).
The third analytic solution (TolVII) was discovered by Tolman Tolman in 1939, and is the case when the mass-energy density $`\rho `$ varies quadratically, that is,
$$\rho =\rho _c[1(r/R)^2].$$
(41)
In fact, this is an adequate representation, as displayed in Fig. 6 for neutron stars more massive than 1.2 M. The equations of state used are listed in Table 1. The largest deviations from this general relation exist for models with extreme softening (GS1, GS2, PCL2) and which have relatively low maximum masses (see Fig. 3). It is significant that all models must, of course, approach this behavior at both extremes $`r0`$ and $`rR`$.
Because the Tolman solution is often overlooked in the literature (for exceptions, see, for example, Refs. Delgaty ; Indians ) it is summarized here. It is useful in establishing interesting and simple relations that are insensitive to the equation of state. In terms of the variable $`x=r^2/R^2`$ and the parameter $`\beta `$, the assumption $`\rho =\rho _c(1x)`$ results in $`\rho _c=15\beta c^2/(8\pi GR^2)`$. The solution of Einstein’s equations for this density distribution is:
$`e^\lambda `$ $`=`$ $`1\beta x(53x),e^\nu =(15\beta /3)\mathrm{cos}^2\varphi ,`$ (42)
$`P`$ $`=`$ $`{\displaystyle \frac{c^4}{4\pi R^2G}}[\sqrt{3\beta e^\lambda }\mathrm{tan}\varphi {\displaystyle \frac{\beta }{2}}(53x),n={\displaystyle \frac{\rho c^2+P}{m_bc^2}}{\displaystyle \frac{\mathrm{cos}\varphi }{\mathrm{cos}\zeta }},`$ (43)
$`\varphi `$ $`=`$ $`(w_1w)/2+\zeta ,\varphi _c=\varphi (x=0),\zeta =\mathrm{tan}^1\sqrt{\beta /[3(12\beta )]},`$ (44)
$`w`$ $`=`$ $`\mathrm{log}[x5/6+\sqrt{e^\lambda /(3\beta )}],w_1=w(x=1).`$ (45)
The central values of $`P/\rho c^2`$ and $`c_s^2`$ are
$$\frac{P}{\rho c^2}|_c=\frac{2}{15}\sqrt{\frac{3}{\beta }}(\frac{c_{sc}}{c})^2,(\frac{c_{sc}}{c})^2=\mathrm{tan}\varphi _c(\mathrm{tan}\varphi _c+\sqrt{\frac{\beta }{3}}).$$
(46)
This solution, like that of Buchdahl’s, is scale-free, with the parameters $`\beta `$ and $`\rho _c`$ (or $`M`$ or $`R`$). Here, $`n`$ is the baryon density, $`m_b`$ is the nucleon mass, and $`c_{sc}`$ is the sound speed at the star’s center. When $`\varphi _0=\pi /2`$, or $`\beta 0.3862`$, $`P_c`$ becomes infinite, and when $`\beta 0.2698`$, $`c_{sc}`$ becomes causal (i.e., $`c`$). Recall that for an incompressible fluid, $`P_c`$ becomes infinite when $`\beta =4/9`$. For the Buchdahl solution, $`P_c`$ becomes infinite when $`\beta =2/5`$ and the causal limit is reached when $`\beta =1/6`$. For comparison, if causality is enforced at high densities, it has been empirically determined that $`\beta <0.34`$ LPMY ; glen .
The general applicability of these exact solutions can be gauged by analyzing the moment of inertia, which, for a star uniformly rotating with angular velocity $`\mathrm{\Omega }`$, is
$$I=(8\pi /3)_0^Rr^4(\rho +P/c^2)e^{(\lambda \nu )/2}(\omega /\mathrm{\Omega })𝑑r.$$
(47)
The metric function $`\omega `$ is a solution of the equation
$$d[r^4e^{(\lambda +\nu )/2}\omega ^{}]/dr+4r^3\omega de^{(\lambda +\nu )/2}/dr=0$$
(48)
with the surface boundary condition
$$\omega _R=\mathrm{\Omega }\frac{R}{3}\omega _R^{}=\mathrm{\Omega }\left[1\frac{2GI}{R^3c^2}\right].$$
(49)
The second equality in the above follows from the definition of $`I`$ and the TOV equation. Writing $`j=\mathrm{exp}[(\nu +\lambda )/2]`$, the TOV equation becomes
$$j^{}=4\pi Gr(P/c^2+\rho )je^\lambda /c^2.$$
(50)
Then, one has
$$I=\frac{2c^2}{3G}\frac{\omega }{\mathrm{\Omega }}r^3𝑑j=\frac{c^2R^4\omega _R^{}}{6G\mathrm{\Omega }}.$$
(51)
Unfortunately, an analytic representation of $`\omega `$ or the moment of inertia for any of the three exact solutions is not available. However, approximations which are valid to within 0.5% are
$`I_{Inc}/MR^2`$ $``$ $`2(10.87\beta 0.3\beta ^2)^1/5,`$ (52)
$`I_{B1}/MR^2`$ $``$ $`(2/34/\pi ^2)(11.81\beta +0.47\beta ^2)^1,`$ (53)
$`I_{TVII}/MR^2`$ $``$ $`2(11.1\beta 0.6\beta ^2)^1/7.`$ (54)
In each case, the small $`\beta `$ limit reduces to the corresponding Newtonian results. Fig. 7 indicates that the Tolman approximation is rather good. Ravenhall & Pethick RP suggested that the expression
$$I_{RP}/MR^20.21/(12u)$$
(55)
was a good approximation for the moment of inertia; however, we find that this expression is not a good overall fit, as shown in Fig. 7. For low-mass stars ($`\beta <0.12`$), none of these approximations is suitable, but it is unlikely that any neutron stars are this rarefied. It should be noted that the Tolman approximation does not adequately approximate some EOSs, especially ones that are relatively soft, such as GM3, GS1, GS2, PAL6 and PCL2.
The binding energy formally represents the energy gained by assembling $`N`$ baryons. If the baryon mass is $`m_b`$, the binding energy is simply $`BE=Nm_bM`$ in mass units. However, the quantity $`m_b`$ has various interpretations in the literature. Some authors assume it is about 940 MeV/$`c^2`$, the same as the neutron or proton mass. Others assume it is about 930 MeV/$`c^2`$, corresponding to the mass of C<sup>12</sup>/12 or Fe<sup>56</sup>/56. The latter would yield the energy released in a supernova explosion, which consists of the energy released by the collapse of a white-dwarf-like iron core, which itself is considerably bound. The difference, 10 MeV per baryon, corresponds to a shift of $`10/9400.01`$ in the value of $`BE/M`$. In any case, the binding energy is directly observable from the detection of neutrinos from a supernova event; indeed, it would be the most precisely determined aspect.
Lattimer & Yahil LY suggested that the binding energy could be approximated as
$$BE1.510^{51}(M/\mathrm{M}_{})^2\mathrm{ergs}=0.084(M/\mathrm{M}_{})^2\mathrm{M}_{}.$$
(56)
This formula, in general, is accurate to about $`\pm 20`$%. The largest deviations are for extremely soft EOSs, as shown in Fig. 8.
However, a more accurate representation of the binding energy is given by
$$BE/M0.6\beta /(10.5\beta ),$$
(57)
which incorporates some radius dependence. Thus, the observation of supernova neutrinos, and the estimate of the total radiated neutrino energy, will yield more accurate information about $`M/R`$ than about $`M`$ alone.
In the cases of the incompressible fluid and the Buchdahl solution, analytic results for the binding energy can be found:
$`BE_{Inc}/M`$ $`=`$ $`{\displaystyle \frac{3}{4\beta }}\left({\displaystyle \frac{\mathrm{sin}^1\sqrt{2\beta }}{\sqrt{2\beta }}}\sqrt{12\beta }\right)1,`$ (58)
$`BE_{B1}/M`$ $`=`$ $`(11.5\beta )\sqrt{12\beta }(1\beta )^11.`$ (59)
The analytic results, the Tolman VII solution, and the fit of Eq. (57) are compared to some recent equations of state in Fig. 9. It can be seen that, except for very soft cases like PS, PCL2, PAL6, GS1 and GS2, both the Tolman VII and Buchdahl solutions are rather realistic.
### 4.3 Crustal Fraction of the Moment of Inertia
In the investigation of pulsar glitches, many models associate the glitch size with the fraction of the moment of inertia which resides in the star’s crust, usually defined to be the region in which dripped neutrons coexist with nuclei. The high-density crust boundary is set by the phase boundary between nuclei and uniform matter, where the pressure is $`P_t`$ and the density $`n_t`$. The low-density boundary is the neutron drip density, or for all practical purposes, simply the star’s surface since the amount of mass between the neutron drip point and the surface is negligible. We define $`\mathrm{\Delta }R`$ to be the distance between the points where the density is $`n_t`$ and zero. One can apply Eq. (47) to determine the moment of inertia of the crust alone with the assumptions that $`P/c^2<<\rho `$, $`m(r)M`$, and $`\omega j\omega _R`$ in the crust. One finds
$$\mathrm{\Delta }I\frac{8\pi }{3}\frac{\omega _R}{\mathrm{\Omega }}_{R\mathrm{\Delta }R}^R\rho r^4e^\lambda 𝑑r\frac{8\pi }{3GM}\frac{\omega _R}{\mathrm{\Omega }}_0^{P_t}r^6𝑑P,$$
(60)
where $`M`$ is the star’s total mass and the TOV equation was used in the last step. In the crust, the fact that the EOS is of the approximate polytropic form $`PK\rho ^{4/3}`$ can be used to find an approximation for the integral $`r^6𝑑P`$, viz.
$$_0^{P_t}r^6𝑑PP_tR^6\left[1+\frac{2P_t}{n_tm_nc^2}\frac{(1+7\beta )(12\beta )}{\beta ^2}\right]^1.$$
(61)
Since the approximation Eq. (57) gives $`I`$ in terms of $`M`$ and $`R`$, and $`\omega _R/\mathrm{\Omega }`$ is given in terms of $`I`$ and $`R`$ in Eq. (49), the quantity $`\mathrm{\Delta }I/I`$ can thus be expressed as a function of $`M`$ and $`R`$ with the only dependence upon the equation of state (EOS) arising from the values of $`P_t`$ and $`n_t`$; there is no explicit dependence upon the higher-density EOS. However, the major dependence is upon the value of $`P_t`$, since $`n_t`$ enters only as a correction. We then find
$$\frac{\mathrm{\Delta }I}{I}\frac{28\pi P_tR^3}{3Mc^2}\frac{(11.67\beta 0.6\beta ^2)}{\beta }\left[1+\frac{2P_t}{n_tm_bc^2}\frac{(1+7\beta )(12\beta )}{\beta ^2}\right]^1.$$
(62)
In general, the EOS parameter $`P_t`$, in the units of MeV fm<sup>-3</sup>, varies over the range $`0.25<P_t<0.65`$ for realistic EOSs. The determination of this parameter requires a calculation of the structure of matter containing nuclei just below nuclear matter density that is consistent with the assumed nuclear matter EOS. Unfortunately, few such calculations have been performed. Like the fiducial pressure at and above nuclear density which appears in the relation Eq. (27), $`P_t`$ should depend sensitively upon the behavior of the symmetry energy near nuclear density.
Choosing $`n_t=0.07`$ fm<sup>-3</sup>, we compare Eq. (62) in Fig. 3 with full structural calculations. The agreement is good. We also note that Ravenhall & Pethick RP developed a different, but nearly equivalent, formula for the quantity $`\mathrm{\Delta }I/I`$ as a function of $`M,R,P_t`$ and $`\mu _t`$, where $`\mu _t`$ is the neutron chemical potential at the core-crust phase boundary. This prediction is also displayed in Fig. 3.
Link, Epstein & Lattimer Link established a lower limit to the radii of neutron stars by using a constraint derived from pulsar glitches. They showed that glitches represent a self-regulating instability for which the star prepares over a waiting time. The angular momentum requirements of glitches in the Vela pulsar indicate that more than $`0.014`$ of the star’s moment of inertia drives these events. If glitches originate in the liquid of the inner crust, this means that $`\mathrm{\Delta }I/I>0.014`$. A minimum radius can be found by combining this constraint with the largest realistic value of $`P_t`$ from any equation of state. Stellar models that are compatible with this constraint must fall to the right of the $`P_t=0.65`$ MeV fm<sup>-3</sup> contour in Fig. 3. This imposes a constraint upon the radius, namely that $`R>3.6+3.9M/\mathrm{M}_{}`$ km.
## 5 The Merger of a Neutron Star with a Low-Mass Black Hole
The general problem of the origin and evolution of systems containing a neutron star and a black hole was first detailed by Lattimer & Schramm LSch , although the original motivation was due to Schramm. Although speculative at the time, Schramm insisted that this would prove to be an interesting topic from the points of view of nucleosynthesis and gamma-ray emission. The contemporaneous discovery HT of the first-known binary system containing twin compact objects, PSR 1913+16, which was also found to have an orbit which would decay because of gravitational radiation within $`10^{10}`$ yr, bolstered his argument. Eventually, this topic formed the core of Lattimer’s thesis thesisL , and the recent spate of activity, a quarter century later, in the investigation of the evolution and mergers of such compact systems has wonderfully demonstrated Schramm’s prescience.
Compact binaries form naturally as the result of evolution of massive stellar binaries. The estimated lower mass limit for supernovae (and neutron star or black hole production) is approximately 8 M. Observationally, the number of binaries formed within a given logarithmic separation is approximately constant, but the relative mass distributions are uncertain. There is some indication that the distribution in binary mass ratios might be flat. The number of possible progenitor systems can then be estimated. Most progenitor systems do not survive the more massive star becoming a supernova. In the absence of a kick velocity it is easily found that the loss of more than half of the mass from the system will unbind it. However, the fact that pulsars are observed to have mean velocites in excess of a few hundred km/s implies that neutron stars are usually produced with large “kick” velocities originating in the supernova explosion. In the case that the kick velocity, which is thought to be randomly directed, opposes the star’s orbital velocity, the chances of the post-supernova binary remaining intact increases. In addition, the separation in a surviving binary will be reduced significantly. Subsequent evolution then progresses to the supernova explosion of the companion. More of these systems survive because in many cases the more massive component explodes. But the surviving systems should both have greatly reduced separations and orbits with high eccentricity.
Gravitational radiation then causes the binary’s orbit to decay, such that circular orbits of two masses $`M_1`$ and $`M_2`$ with initial semimajor axes $`a`$ satisfying
$$a<2.8[M_1M_2(M_1+M_2)/\mathrm{M}_{}^3]^{1/4}\mathrm{R}_{},$$
(63)
will fully decay within the age of the Universe ($`10^{10}`$ yr). Highly eccentric orbits will decay much faster, as shown in Fig. 10. The dashed curve shows the inverse of the factor Peters by which the gravitational wave luminosity of an eccentric system exceeds that of a circular system:
$$f=(1+73e^2/24+37e^4/96)(1e^2)^{7/2}.$$
(64)
Because the eccentricity also decays, the exact reduction factor is not as strong as $`1/f`$. A reasonable approximation to the exact result is $`f^{3/4}`$, shown by the dotted line in Fig. 10. The coefficient 2.8 in Eq. (63) is increased by a factor of $`f^{3/16}`$ or about 2 for moderate eccentricities.
Ref. LSch argued that mergers of neutron stars and black holes, and the subsequent ejection of a few percent of the neutron star’s mass, could easily account for all the r-process nuclei in the cosmos. Ref. LSch is also the earliest reference to the idea that compact object binary mergers are associated with gamma-ray bursts. A later seminal contribution by Eichler, Livio, Piran & Schramm eichler argued that mergers of neutron stars occur frequently enough to explain the origin of gamma-ray bursters.
Since the timescale of gamma-ray bursts, being of order seconds to several minutes, is much longer than the coalescence timescale of a binary merger (which is of order the orbital frequency at the last stable orbit, a few milliseconds), it is believed that a coalescence involves the formation of an accretion disc. Although neutrino emission from accreting material, resulting in neutrino-antineutrino annihilation along the rotational axis, has been proposed as a source of gamma rays, it seems more likely that amplification of magnetic fields within the disc might trigger observed bursts. In either case, the lifetime of the accretion disc is still problematic, if it is formed by the breakup of the neutron star near the Roche limit. Its lifetime would probably be only about a hundred times greater than the orbital frequency, or less than a second. However, this timescale would be considerably enhanced if the accretion disc could be formed at larger radii than the Roche limit. A possible mechanism is stable mass transfer from the neutron star to the black hole that would cause the neutron star to spiral away as it loses mass Kochanek ; PZ .
The classical Roche limit is based upon an incompressible fluid of density $`\rho `$ and mass $`M_2`$ in orbit about a mass $`M_1`$. In Newtonian gravity, this limit is
$$R_{Roche,Newt}=(M_1/0.0901\pi \rho )^{1/3}=19.2(M_1/\mathrm{M}_{}\rho _{15})^{1/3}\mathrm{km},$$
(65)
where $`\rho _{15}=\rho /10^{15}`$ g cm<sup>-3</sup>. Using general relativity, Fishbone Fishbone found that at the last stable circular orbit (including the case when the black hole is rotating) the number 0.0901 in Eq. (65) becomes 0.0664. In geometrized units, $`R_{Roche}/M_1=13(14.4)(M_1^2\rho _{15}/\mathrm{M}_{}^2)^{1/3}`$, where the numerical coefficient refers to the Newtonian (last stable orbit in GR) case. In other words, if the neutron star’s mean density is $`\rho _{15}=1`$, the Roche limit is encountered beyond the last stable orbit if the black hole mass is less than about 5.9 M. Thus, for small enough black holes, mass overflow and transfer from the neutron star to the black hole could begin outside the last stable circular orbit. And, as now discussed, the mass transfer may proceed stably for some considerable time. In fact, the neutron star might move to 2–3 times the orbital radius where mass transfer began. This would provide a natural way to lengthen the lifetime of an accretion disc, by simply increasing its size.
The final evolution of a compact binary is now discussed. Define $`q=m_{ns}/M_{BH}`$, $`\mu =m_{ns}M_{BH}/M`$, and $`M=M_{BH}+m_{ns}`$, where $`m_{ns}`$ and $`M_{BH}`$ are the neutron star and black hole masses, respectively. The orbital angular momentum is
$$J^2=G\mu ^2Ma=GM^3aq^2/(1+q)^4.$$
(66)
We can employ Paczyński’s pacz formula for the Roche radius of the secondary:
$$R_{\mathrm{}}/a=0.46[q/(1+q)]^{1/3},$$
(67)
or a better fit by Eggleton eggleton :
$$R_{\mathrm{}}/a=0.49[.6+q^{2/3}\mathrm{ln}(1+q^{1/3})]^1.$$
(68)
The orbital separation $`a`$ at the moment of mass transfer is obtained by setting $`R_{\mathrm{}}=R`$, the neutron star radius. For stable mass transfer, the star’s radius has to increase more quickly than the Roche radius as mass is transferred. Thus, we must have, using Paczyński’s formula,
$$\frac{d\mathrm{ln}R}{d\mathrm{ln}m_{ns}}\alpha \frac{d\mathrm{ln}R_{\mathrm{}}}{d\mathrm{ln}m_{ns}}=\frac{d\mathrm{ln}a}{d\mathrm{ln}m_{ns}}+\frac{1}{3}$$
(69)
for stable mass transfer. $`\alpha `$ is defined in this expression, and is shown in Fig. 11 for a typical EOS. If the mass transfer is conservative, than $`\dot{J}=\dot{J}_{GW}`$, where
$$\dot{J}_{GW}=\frac{32}{5}\frac{G^{7/2}}{c^5}\frac{\mu ^2M^{5/2}}{a^{7/2}}=\frac{32}{5}\frac{G^{7/2}}{c^5}\frac{q^2M^{9/2}}{(1+q)^4a^{7/2}}$$
(70)
and
$$\frac{\dot{J}}{J}=\frac{\dot{a}}{2a}+\frac{\dot{q}(1q)}{q(1+q)}.$$
(71)
This leads to
$$\dot{q}\left(\frac{\alpha }{2}+\frac{5}{6}q\right)\frac{32}{5}\frac{G^3}{c^5}\frac{q^2M^3}{(1+q)a^4}.$$
(72)
Since $`m_{ns}<M_{BH}`$, $`\dot{q}0`$, and the condition for stable mass transfer is simply $`q5/6+\alpha /2`$. For moderate mass neutron stars, $`\alpha 0`$, so in this case the condition is simply $`q5/6`$, which might even be achievable in a binary neutron star system. Had we used the more exact formula of Eggleton, Eq. (68), we would have found $`q0.78`$. Note that it has often been assumed that $`Rm_{ns}^{1/3}`$ in such discussions PZ , which is equivalent to $`\alpha =1/3`$. This is unjustified, and results in the upper limit $`q=2/3`$ which might inappropriately rule out stable mass transfer in the case of two neutron stars.
A number of other conditions must hold for stable mass transfer to occur. First, the orbital separation $`a`$ at the onset must exceed the last stable orbit around the black hole, so that $`a>6GM_{BH}/c^2`$, or
$$q6\frac{R_{\mathrm{}}}{a}\frac{GM_{BH}}{Rc^2}.$$
(73)
Second, the tidal bulge raised on the neutron star must stay outside of the black hole’s Schwarzschild radius. Kochanek Kochanek gives an estimate of the height of the tidal bulge needed to achieve the required mass loss rate:
$$\frac{\mathrm{\Delta }r}{R}=\left[\frac{\dot{q}}{\beta _t(1+q)\mathrm{\Omega }}\right]^{1/3},$$
(74)
where $`\beta _t`$ is a dimensionless parameter of order 1 and $`\mathrm{\Omega }=G^{1/2}M^{1/2}/a^{3/2}`$ is the orbital frequency. For $`\dot{q}`$ we use the equality in Eq. (72), which is equivalent to
$$R_{sh}=2GM_{BH}/c^2aR\mathrm{\Delta }r.$$
(75)
Finally, so that the assumption of a Roche geometry is valid, it should be possible for tidal synchronization of the neutron star to be maintained. Bildstein & Cutler BC considered this, and derived an upper limit for the separation $`a_{syn}`$ at which tidal synchronization could occur by integrating the maximum torque on the neutron star as it spirals in from infinity and finding where the neutron star spin frequency could first equal the orbital frequency. They find
$$a_{syn}\frac{M_{BH}^2m_{ns}^2}{400M^3}\left(\frac{R}{m_{ns}}\right)^6,$$
(76)
which translates to
$$400\left(\frac{GM_{BH}}{Rc^2}\right)^5\frac{a}{R_{\mathrm{}}}\frac{(1+q)^3}{q}1.$$
(77)
Next we consider the effect of putting some of the angular momentum into an accretion disc. Following the discussion of Ref. BC , we assume an accretion disc contains an amount of angular momentum that grows at the rate
$$\dot{J}_d=(1f)M^{3/2}a^{1/2}(1+q)^4\dot{q},$$
(78)
where $`f`$ is a parameter, taken to be a fit to the numerical results of Hut & Paczyński HP :
$$f=5q^{1/3}/33q^{2/3}/2.$$
(79)
We then find the new condition for angular momentum conservation to be
$$\dot{J}+\dot{J}_d=\dot{J}_{GW},$$
(80)
which yields
$$\dot{q}\left[\frac{\alpha }{2}\frac{1}{6}+\frac{fq^2}{1+q}\right]\frac{32}{5}\frac{G^3}{c^5}\frac{q^2M^3}{(1+q)a^4}.$$
(81)
Therefore, the new condition for stable mass transfer is
$$(q^2f)/(1+q)\alpha /21/6.$$
(82)
The case $`f=1`$ corresponds to neglecting the existance of an accretion disc.
It remains to determine when an accretion disc is likely to form. Initially, matter flowing from the neutron star to the black hole through the inner Lagrangian point passes close to the black hole and falls in. However, as the neutron star spirals away, the accretion stream trajectory moves outside the Schwarzschild radius. When the trajectory doesn’t even penetrate the marginally stable orbit, an accretion disc will begin to form. Particle trajectory computations of the Roche geometry by Shore, Livio & van den Huevel SLv suggest that its closest approach to the black hole is
$$R_c=a(1+q)(0.50.227\mathrm{ln}q)^4.$$
(83)
Equating $`R_c`$ to $`6GM_{BH}/c^2`$ yields
$$(0.50.227\mathrm{ln}q)^4(1+q)6\frac{GM_{BH}}{Rc^2}\frac{R_{\mathrm{}}}{a}.$$
(84)
These constraints and allowed regions for stable mass transfer are shown in Fig. 12. Apparently, stable mass transfer ceases when $`m_{ns}0.14`$ M if the formation of an accretion disc is ignored. If the effects of disc formation are included, the stable mass transfer ceases when $`m_{ns}0.22`$ M. In both cases, the neutron star mass remains above its minimum mass (about 0.09 M for the equation of state used here). Thus, the neutron star does not “explode” by reaching its minimum mass.
Fig. 13 shows the time development of the orbital separation $`a`$ and the neutron star’s mass and radius during the inspiral and stable mass transfer phases. Solid lines are calculated assuming there is no accretion disc formed, while dashed lines show the effects of accretion disc formation. The time evolutions during stable mass transfer are obtained from Eq. (72) and Eq. (81), using $`\dot{m}_{ns}=\dot{q}M/(1+q)^2`$. With disc formation, the mass transfer is accelerated and the duration of the stable mass transfer phase is shortened considerably. Also, the neutron star spirals out to a smaller radius, and does not lose as much mass, as in the case when the accretion disc is ignored.
Therefore, if stable mass transfer can take place, the timescale over which mass transfer occurs will be much longer than an orbital period, and lasts perhaps a few tenths of a second. This is not long enough to explain gamma-ray bursts. However, we have also seen the likelihood that an accretion disc forms is quite large. Furthermore, the accretion disc extends to about 100 km. Even though this is considerably less than Ref. PZ estimated, the lifetime of such an extended disc is considerable. To order of magnitude, it is given by the viscous dissipation time, or
$$\tau _{visc}\frac{D^2}{\alpha c_sH}.$$
(85)
Here $`D`$ is the radial size of the disc, $`\alpha `$ is the disc’s viscosity parameter, $`c_s`$ is the sound speed and $`H`$ is the disc’s thickness. Note that $`c_s\mathrm{\Omega }H`$ where $`\mathrm{\Omega }=2\pi /P=\sqrt{GM_{BH}/D^3}`$ is the Kepler frequency. Thus,
$$\tau _{visc}\frac{P}{2\pi \alpha }\left(\frac{D}{H}\right)^2.$$
(86)
Since the magnitude of $`\alpha `$ is still undetermined, and usually quoted Bran to be about 0.01, and $`H`$ is likely to be of order $`R`$, we find $`\tau _{visc}230`$ s for our case. This alleviates the timescale problem for these models. Numerical simulations of such events are in progress, and it remains to be seen if a viable gamma-ray burst model from neutron star–black hole coalescence is possible. If it is, a great deal of the credit should rest with Dave.
We thank Ralph Wijers for discussions concerning accretion disks. |
warning/0002/astro-ph0002261.html | ar5iv | text | # Cooling curves and initial models for low–mass white dwarfs (<0.25𝑀_⊙) with helium core
## 1 Introduction
Kippenhahn, Kohl & Weigert (1967) were the first who followed the formation of helium white dwarfs (WD) of low mass in a binary system. The evolution of a helium WD of 0.26$`M_{}`$ (remnant) was investigated by Kippenhahn, Thomas & Weigert (1968) who found that a hydrogen flash can be initiated near the base of the hydrogen rich envelope. The energy of the flash is sufficient to cause the envelope to expand to giant dimensions and hence it may be possible that another short term Roche lobe filling can occur.
In Webbink (1975), models of a helium white dwarf were constructed by formally evolving a model from the homogeneous zero–age main sequence with the reduction of the mass of the hydrogen–rich envelope. When the mass of the envelope is less than some critical value, the model contracts adopting white dwarf dimensions. Webbink found that thermal flashes do not occur for WDs less massive than 0.2$`M_{}`$. Alberts et al. (1996) have confirmed Webbink’s finding that low–mass white dwarfs do not show thermal flashes and the cooling age for WDs of mass $`M_{wd}`$$``$0.20$`M_{}`$ can be considerably underestimated if using the traditional WD cooling curves which were constructed for $`M_{wd}`$$`>`$0.3$`M_{}`$ (Iben & Tutukov 1986, IT 86).
Recently, Hansen & Phinney (1998a – HP98) and Benvenuto & Althaus (1998 – BA98) investigated the effect of different mass of the hydrogen layer ($`10^8\mathrm{M}_{\mathrm{env}}/\mathrm{M}_{}4\times 10^3`$) on the cooling evolution of $`0.15\mathrm{M}_{\mathrm{He}}/\mathrm{M}_{}0.5`$ helium WDs. In both calculations (BA98 and HP98) the mass of the hydrogen envelope left on the top of white dwarf has been taken as free parameter. BA98 found that thick envelopes appreciably modify the radii and surface gravities of no–H models, especially in the case of low–mass helium white dwarfs.
Driebe et al. (1998 – DSBH98) present a grid of evolutionary tracks for low-mass white dwarfs with helium cores in the mass range from 0.179 to 0.414$`M_{}`$. The tracks are based on a 1$`M_{}`$ model sequence extending from the pre–main sequence stage up to the tip of red giant branch. Applying large mass loss rates forced the models to move off the giant branch and evolve across the Hertzsprung–Russell diagram and down the cooling branch. They found that hydrogen flashes take place only for two model sequences, 0.234$`M_{}`$ and 0.259$`M`$, and for very low–mass WDs the hydrogen shell burning remains dominant even down to effective temperatures well below 10 000 K. According to our previous calculations (Ergma, Sarna & Antipova, 1998) we find that for a low–mass white dwarf with a helium core, which was formed during low–mass binary evolution (after detachment from the Roche lobe), the hydrogen layer left on the top of the helium core is much thicker ($`16\times 10^2\mathrm{M}_{}`$ with $`\mathrm{X}_{\mathrm{surf}}`$ ranging from 0.3 to 0.52) than used in cooling calculation by HP98 and BA98. Also in DSBH98 (see their Table 1), for the two lowest total remnant masses the envelope mass value is smaller that obtained in our calculations.
## 2 The main aim
Low–mass helium white dwarfs are present in millisecond binary pulsars and double degenerate systems. This gives a unique opportunity to test the cooling age of the WD in a binary and, especially in the case of millisecond binary pulsars, allows for age determinations for neutron stars that are independent of their rotational history.
## 3 The evolutionary code
The evolutionary sequences we have calculated are comprised of three main phases:
$``$ detached evolution lasting until the companion fills its Roche lobe on the time–scale $`t_d`$;
$``$ semi–detached evolution (non–conservative in our calculations) on the time–scale $`t_{sd}`$; $`t_0`$=$`t_d+t_{sd}`$;
$``$ a cooling phase of the WD on the time–scale $`t_{cool}`$ (the final phase during which a system with a ms pulsar + low–mass helium WD is left behind). The total evolutionary time is $`t_{evol}=t_0+t_{cool}`$.
The duration of the detached phase is somewhat uncertain; it may be determined either by the nuclear time–scale or by the much shorter time–scale of the orbital angular momentum loss owing to the magnetized stellar wind.
In our calculations we assume that the semi–detached evolution of a binary system is non–conservative, i.e. the total mass and angular momentum of the system are not conserved. We can express the total orbital angular momentum (J) of a binary system as
$$\frac{\dot{J}}{J}=\frac{\dot{J}}{J}|_{SML}+\frac{\dot{J}}{J}|_{MSW}+\frac{\dot{J}}{J}|_{GR},$$
(1)
where the terms on the right hand side are due to: stellar mass angular momentum loss from the system, magnetic stellar wind braking, and gravitational wave radiation.
### 3.1 Stellar mass angular momentum loss
The formalism which we have adopted is described in Muslimov & Sarna (1993). We introduce the parameter $`f_1`$ characterizing the loss of mass from the binary system and defined by the relations,
$$\dot{M}=\dot{M}_2f_1and\dot{M}_1=\dot{M}_2(1f_1),$$
(2)
where $`\dot{M}`$ is the mass–loss rate from the system, $`\dot{M}_2`$ is the rate of mass loss from the donor (secondary) star and $`\dot{M}_1`$ is the accretion rate onto the neutron star (primary). The matter leaving the system will carry off its intrinsic angular momentum in agreement with formula
$$\frac{\dot{J}}{J}|_{SML}=f_1f_2\frac{M_1\dot{M}_2}{M_2M}yr^1,$$
(3)
where $`M_1`$ and $`M_2`$ are the masses of the neutron star and donor star, respectively and M=$`M_1`$+$`M_2`$. Here we have introduced the additional parameter $`f_2`$, which describes the efficiency of the orbital angular momentum loss from the system due to a stellar wind (Tout & Hall 1991). In our calculations we have $`f_2`$=1 and $`f_1`$=1; we calculate the fully non–conservative case, although additional calculations with $`f_1`$ = 0.9 and 0.5 (with $`f_2`$=1) give similar results. A similar result to ours was found by Tauris (1996), who showed that the change in orbital separation due to mass transfer in LMXB (low–mass X-ray binaries) as a function of the fraction of exchanged matter $`f_1`$ which is lost from system is small (for 0.5$`f_11`$). To understand whether the system evolution is conservative or non–conservative is not easy in the case of a rapidly rotating neutron star; no easy solution can be found. We propose as one possibility a factor which may help us to distinguish between the two cases – the surface magnetic field of the neutron star and its evolution during the accretion.
### 3.2 Magnetic stellar wind braking
We also assume that the donor star, possessing a convective envelope, experiences magnetic braking (Mestel 1968; Mestel & Spruit 1987; Muslimov & Sarna 1995), and, as a consequence, the system loses its orbital angular momentum. For a magnetic stellar wind we used the formula for the orbital angular momentum loss
$$\frac{\dot{J}}{J}|_{MSW}=3\times 10^7\frac{M^2R_2^2}{M_1M_2a^5}yr^1,$$
(4)
where $`a`$ and $`R_2`$ are the separation of the components and the radius of the donor star in solar units.
### 3.3 Gravitational wave radiation
For systems with very short orbital periods, during the final stages of their evolution we also take into account the loss of orbital angular momentum due to emission of gravitational radiation (Landau & Lifshitz 1971):
$$\frac{\dot{J}}{J}|_{GR}=8.5\times 10^{10}\frac{M_1M_2M}{a^4}yr^1$$
(5)
The mass and accompanying orbital angular momentum loss from these system are poorly understood problems in the evolution of binary stars. As is well known, the variation of the angular momentum depends critically on the assumed model (Ergma et al. 1998). In the case of binary systems with ms pulsar typically two different models concerning the mass ejection and angular momentum loss can be adopted. The first is that the amount of angular momentum lost per 1 gram of ejected matter is equal to the average orbital angular momentum of 1 gram of the binary. The second is that the matter that flows from the companion star onto the neutron star (after accretion) is ejected isotropically with the specific angular momentum of the neutron star. In this paper, for our non–conservative approach we have adopted the first model. This affects significantly our results on the semi–detached evolution (see fig. 2 in Ergma et al. 1998), but very little changes the cooling time–scale of the helium white dwarf.
### 3.4 Illumination of the donor star
In all cases we have included the effect of illumination of the donor star by the millisecond pulsar. In our calculations we assume that illumination of the component by the hard (X–ray and $`\gamma `$–ray) radiation from the millisecond pulsar leads to additional heating of its photosphere (Muslimov & Sarna 1993). The effective temperature $`T_{eff}`$ of the companion during the illumination stage is determined from the relation
$$L_{in}+P_{ill}=4\pi \sigma R_2^2T_{eff}^4,$$
(6)
where $`L_{in}`$ is the intrinsic luminosity corresponding to the radiation flux coming from the stellar interior and $`\sigma `$ is the Stefan–Boltzmann constant.
$`P_{ill}`$ is the millisecond pulsar radiation that heats the photosphere, which is determined by
$$P_{ill}=f_3\left(\frac{R_2}{2a}\right)^2L_{rot}$$
(7)
and $`L_{rot}`$ is “rotational luminosity” of the neutron star due to magneto–dipole radiation (plus a wind of relativistic particles)
$$L_{rot}=\frac{2}{3c^3}B^2R_{ns}^6\left(\frac{2\pi }{P_p}\right)^4,$$
(8)
where $`R_{ns}`$ is the neutron star radius, B is the value of the magnetic field strength at the neutron star and $`P_p`$ is the pulsar period. $`f_3`$ is the factor characterizing the efficiency of transformation of irradiation flux into thermal energy (in our case we take $`f_3=2\times 10^3`$). Note that in our calculations the effect of irradiation is formally treated by means of modification of the outer boundary condition, according to relation (6).
In this paper we do not follow the magnetic field and pulsar period ($`P_p`$) evolution, as we did in our earlier papers (Muslimov & Sarna 1993, Ergma & Sarna 1996). We were mainly interested in finding initial models for low–mass helium white dwarfs and in investigating the initial cooling phase of these low–mass helium white dwarfs. From earlier calculations we know that if the magnetic field strength is greater than about $`10^9`$G, the neutron star spins–up to tenths and hundreds of milliseconds, rather than several milliseconds. This leads to a situation where the pressure of the magneto–dipole radiation is insufficient to eject matter from the system. Also from our previous calculations (see for example Ergma & Sarna 1996) we find that after accretion of a maximum of about 0.2$`M_{}`$, the neutron star has spun–up to millisecond periods if B$`<10^9`$G. Therefore in this paper we accept that after accretion of 0.2$`M_{}`$ the neutron star spins–up to about 2 ms. After spin–up the pulsar irradiation is strong enough to prevent accretion, and at this moment we include non–conservative mass loss from the system as described above.
During the initial high mass accretion phase ($`\dot{M}_210^810^9M_{}yr^1`$, $`t_{acc}10^710^8`$yrs) the system may be observed as a bright low–mass X–ray binary (LMXB). It is necessary to point out that majority of LMXBs for which orbital period determinations are available (21 systems out of 24 according to van Paradijs catalogue 1995), have orbital period of less than one day. These systems therefore cannot be the progenitors of the majority of low–mass helium white dwarf + millisecond pulsar binary systems. A lack of LMXB systems with orbital period between 1 – 3 days does not allow us to make a direct comparison between the observational data and the results of our calculations.
### 3.5 The code
The models of the stars filling their Roche lobes were computed using a standard stellar evolution code based on the Henyey–type code of Paczyński (1970), which has been adapted to low–mass stars. The Henyey method involves iteratively improving a trail solution for the whole star. During each iteration, corrections to all variables at all mesh points in the star are evaluated using the Newton–Raphson method for linearised algebraic equations (see for example Hansen & Kawaler 1994). The Henyey method extended to calculate stellar evolution with mass loss, as adopted here, is well explain by Ziółkowski (1970). We note here that our code makes use of the stationary envelope technique, which was developed early on in the life of our code in order to save disc space (Paczyński 1969). This method makes the assumption that the surface 0.5 – 5% (by mass) of the star is not significantly affected by nuclear processes, such that it can be treated to a good approximation as homogeneous region (in composition) throughout the whole evolutionary calculation. During the cooling phase we assume that the static envelope is the surface 0.5% of the star. This assumption is valid during the flashes because the time–scale is longer than thermal time–scale of the envelope. We tested the possibility that the algorithm for redistributing meshpoints introduces numerical diffusion into the composition profile. We find that if such numerical diffusion is real, it has only a marginal influence on the hydrogen profile. We would also like to note that in the heat equation we neglect the derivative with respect to molecular weight, since its effect is small. Convection is treated with the mixing–length algorithm proposed by Paczyński (1969). We solve the problem of radiative transport by employing the opacity tables of Iglesias & Rogers (1996). Where the Iglesias & Rogers (1996) tables are incomplete, we have filled the gaps using the opacity tables of Huebner et al. (1977). For temperatures less than 6000 K we use the opacities given by Alexander & Ferguston (1994) and Alexander (private communication). The contribution from conduction present in the opacity tables of Huebner et al. (1977) has been included by us in the other tables, since they don’t include it (Haensel, private communication). The equation of state (EOS) includes radiation and gas pressure, which is composed of the ion and electron pressure. Contribution to the EOS owing to the non–ideal effects of Coulomb interaction and pressure ionization which influence the EOS, as discussed by Pols et al. (1995), have not been included in our program, and for this reason we stopped our cooling calculations before these effects become important. During the initial phase of cooling, the physical conditions in the hot white dwarfs are such that these effects are usually small.
## 4 Evolutionary calculations
We perform our evolutionary calculations for binary systems initially consisting of a 1.4 $`\mathrm{M}_{}`$ neutron star (NS) and a slightly evolved companion (subgiant) of two masses, 1 and 1.5, and four chemical compositions (Z: 0.003, 0.01, 0.02, 0.03). We have produced (Table 1) a number of evolutionary tracks corresponding to the different possible values of the initial orbital period (ranging from 0.7 to 3.0 days) at the beginning of mass transfer phase.
## 5 The Results
In Table 1 we list the characteristic of the cooling phase of the WD, $`\mathrm{t}_{\mathrm{cool}}`$, and the maximum possible evolution time of a system, $`\mathrm{t}_{\mathrm{evol}}`$, which is a sum of times of detached (determined by nuclear evolution), semi–detached, and cooling phases. The cooling is the last phase of evolution of the WD, and in our calculations starts at the end of RLOF. The cooling time, $`\mathrm{t}_{\mathrm{cool}}`$, is limited to an initial cooling stage during which the WD cools until its central temperature has decreased by 50 $`\%`$ of its maximum value. From Table 1 it is clearly seen that to produce short orbital period systems in a time–scale shorter than Hubble time it is necessary either to have low Z or a more massive secondary.
In our calculations the donor star fills its Roche lobe while it is evolving through the Hertzsprung gap, and therefore it transfers mass on its companion in a thermal time–scale.
Figure 1 show the evolutionary cooling sequences for models 20 and 22 (more details in Table 1). Model 20 presents the case with stable hydrogen burning. Model 22 shows the case when the thermal instability of the hydrogen–burning shell occurs. The first flash is not strong enough to allow the star to overflow its Roche lobe, but during the second flash the radius of the secondary increases to fill its Roche lobe and short–time Roche lobe overflow (RLOF) occurs.
In Table 2 we present the mass–radius relationship for WDs from our calculations, DSBH98, the Wood models, and the Hamada & Salpeter (1961) zero–temperature helium WD models calculated for a surface temperature of 8500 K (as in van Kerkwijk, Bergeron & Kulkarini 1996 for PSR1012+5307). Comparison of the numbers demonstrate that for WD masses of $`<`$ 0.25$`M_{}`$, the results of our calculations differ significantly from a simple extrapolation obtained from the cooling curves (Wood 1990) performed for carbon WDs with the thick hydrogen envelopes. In addition comparing the cooling time–scales of HP98 and BA98 with those of Webbink and our models, shows differences of an order of magnitude (Table 3) for WD masses of $`<`$ 0.25$`M_{}`$.
## 6 Hydrogen flash burning
The problem of unstable hydrogen shell burning in low–mass helium WDs was first discussed in the literature more than 30 years ago (Kippenhahn, Thomas & Weigert 1968). Recently, Alberts et al. (1996) have claimed that they do not see any thermal flashes that result from thermally unstable shell–burning, as reported in papers IT86 and Kippenhahn, Thomas & Weigert (1968). Webbink (1975) found that in none of his model sequences, such a severe thermal runaway as described by Kippenhahn et al. (1968) was found, although mild flashes for M$`>`$0.2 $`M_{}`$ did take place. Alberts et al. found that even reducing the time step to 50–100 years would not lead to thermally unstable shell–burning for $`M_{wd}`$$`<`$0.25 $`M_{}`$. In DSBH98, thermal instabilities of the hydrogen–burning shell occurs in their two models, 0.234$`M_{}`$ and 0.259$`M_{}`$. They concluded that hydrogen flashes take place only in the mass interval 0.21$``$ $`M/M_{}0.3`$.
According to our computations, low–mass helium WDs with masses more than 0.183$`M_{}`$ (Z=0.03), 0.192$`M_{}`$ (Z=0.02), 0.198$`M_{}`$ (Z=0.01) and 0.213$`M_{}`$, (Z=0.003) may experience up to several hydrogen flashes before they enter the cooling stage. In Table 4 we present several characteristics for the computed flashes. We discussion two kinds of flashes: in the first case (in Table 4 shown as “1”), during the flash the secondary does not fill its Roche lobe i.e. the mass of the white dwarf does not change, and in the second case (“2”), during the unstable hydrogen burning phase the secondary fills its Roche lobe and the system again enters into a very short duration accretion phase (see Table 4). We introduce four time–scales to describe the flash behaviour: (i) the flash rise time–scale $`\mathrm{\Delta }t_1`$, which is the time for the luminosity to increase from minimum to maximum value (typically this value is between few $`\times 10^6`$ to few $`\times 10^7`$ yrs – third column in the Table 4); (ii) the flash decay time–scale $`\mathrm{\Delta }t_2`$, which is the time for the luminosity to decrease to the initial value (typically from few hundred thousand to few tenth million years); (iii) $`\mathrm{\Delta }`$ T is the recurrence time between two successive flashes (iv) $`\mathrm{\Delta }t_{acc}`$ is the duration of the accretion phase when the secondary fills its Roche lobe during hydrogen shell flash.
For all sequences with several unstable hydrogen shell burning stages (usually for case “1”), the first flash is the weakest. In the majority of cases when the flash forces the star to fill its Roche lobe, only one flash takes place. For four cases we found two successive flashes with Roche lobe overflow (models 17, 23, 24, 31), and for another two cases (models 47, 53) to the first flash is not powerful enough to force the secondary fill its Roche lobe, but during the second flash it is.
How does the hydrogen flash burning influence the cooling time–scale? In Fig.2, the luminosity and nuclear energy production rates versus cooling time for models 20 and 22 are shown. Model 20 shows stationary hydrogen burning and model 22, hydrogen flash burning. Although before flash model 22 was more luminous than model 20, later the situation is reversed. After the flash, the burning mass of the hydrogen rich envelope in model 22 has decreased to 0.0116$`M_{}`$, whereas the mass of the hydrogen envelope in model 20, in which stationary hydrogen burning occurs, is almost twice as large (0.0241$`M_{}`$). If we look at how the maximum nuclear energy rate behaves with cooling time, we can see that after the flash in model 22, the maximum energy production rate is less than in model 20 (stationary hydrogen burning).
In Fig. 3 we present the behaviour of log$`T_{eff}`$, log $`ϵ_{nuc}`$ and log L/$`L_{}`$, and in Fig.4 log $`R_{wd}`$, $`M_{env}`$ and $`M_f/M_{}`$ as a function of cooling time for model 7. Before the helium white dwarf enters the final cooling phase, four unstable hydrogen flash burnings occur. The same parameters for model 17 (with RLOF) are shown in Figs. 5 and 6.
To investigate in more detail how the flashes develop, we show in Fig. 7 the evolution of the white dwarf radius (upper panel), nuclear energy generation rate (upper middle panel), maximum shell temperature and central temperature (lower middle panel) and the surface luminosity (lower panel) as a function of computed model number. In Fig. 7, as vertical dashed lines we marked several time–scales which characterize the flash behaviour (for numbers see Table 4). $`\mathrm{\Delta }t_1`$ and $`\mathrm{\Delta }t_2`$ describe the rise and decay times; the first characterizes the nuclear shell burning time–scale ($`\tau _{nuc}^{shell}`$), the second the Kelvin–Helmholtz (thermal) envelope time–scale modified by nuclear shell burning ($`\mathrm{\Delta }t_2=\sqrt{\tau _{KH}^{env}\tau _{nuc}^{shell}}`$). The accretion time ($`\mathrm{\Delta }t_{acc}`$) is described by the square of the Kelvin–Helmholtz time–scale. The radiative diffusion time is defined as the Kelvin–Helmholtz time–scale of the extended envelope above the shell ($`\mathrm{\Delta }t_{rd}=\tau _{KH}^{env}`$). The shape of the first flash on Fig. 7 shows some characteristic changes which are connected with physical processes in the stellar interior. At the beginning of the flash the luminosity increases due to the more effective hydrogen burning in the shell source. After reaching a local maximum, the luminosity then decreases while the nuclear energy generation rate is still increasing rapidly. This decrease of the surface luminosity is due to a temperature inversing forming below the hydrogen shell. The energy generated in the hydrogen shell splits into two fluxes; coming outwards and going inwards. The helium core is heated effectively by the shell nuclear source – the central temperature increases by 2%. On Fig. 8 the evolution of the luminosity and temperature profiles during the $`\mathrm{\Delta }t_1`$ and $`\mathrm{\Delta }t_2`$ phases are shown. We clearly see how the inversion profile evolves and how the luminosity wave moves into the surface.
The nuclear energy generation rate in the shell has a maximum value far away from maximum surface luminosity. This is because the luminosity front is moving towards the stellar surface in a time–scale described by radiative diffusion ($`\mathrm{\Delta }t_{rd}`$). After reaching a maximum value, the luminosity starts to decrease and the energy generation rate also declines in the hydrogen shell over a time-scale $`\mathrm{\Delta }t_2`$ (for a contracting envelope) the luminosity decreases to the minimum value. During the first flash, the stellar radius does not fill the inner Roche lobe. In the second and third flashes we have short episodes of super–Eddington mass transfer (see Table 4). During the RLOF phase, the orbital period slightly increases and the subgiant companion evolves quickly from spectral type F0 to A0.
As already pointed out, for several cases the secondary fills its Roche lobe and the system enters an accretion phase. During RLOF, the mass accretion rate is about three orders of magnitude greater than the Eddington limit (Fig. 9). All the accreted matter will be lost from the system ($`\mathrm{\Delta }M_{acc}0.00010.001M_{}`$). The accretion phase is very short, usually less than 1000 years (ranging from 160 to 2500 yrs – see Table 4). During the short super–Eddington accretion phase the system is a very bright X–ray source, with orbital period between 2 to 8 days.
We notice that during the flash the evolutionary time step strongly decreases and may be as short as several years.
## 7 Role of binarity in the cooling history of the low–mass white dwarf
DSBH98 modelled single star evolution and produced white dwarfs with various masses by applying large mass loss rates at appropriate positions in the red–giant branch to force the models to move off the giant branch. To show how binarity influences the final fate of the white dwarf cooling, we have computed extra sequences (1.0+ 1.4 $`M_{}`$, Z=0.02, $`P_i`$=2.0 days) where we did not take into account that the star is in a binary system e.g. during hydrogen shell flash we do not allow RLOF. In complete binary model calculation, only one shell flash occurs accompanied with RLOF, whereas for the single star model calculation, four hydrogen shell flashes take place. Due to RLOF, the duration of the flash phase is 2.7$`\times 10^6`$yrs; if we do not include binarity the duration of the flash phase is 1.8$`\times 10^8`$yrs. However, the cooling time for helium white dwarfs less massive than 0.2$`M_{}`$ is not significantly changed. This is because the duration of flash phase is very short in comparison to the normal cooling phase (towards the white dwarf region). However, the effect of binarity will be important for the cooling history of more massive helium white dwarfs. In Fig. 10 both cases of evolution on the Hertzsprung–Russell diagram are shown – on the left panel Roche lobe overflow is not allowed, on right panel RLOF takes place.
## 8 Application to individual systems
Below we discuss the observational data for several systems for which results of our calculations may be applied, by taking into account the orbital parameters of the system, the pulsar spin–down time, and the white dwarf cooling timescale.
### 8.1 PSR J0437–4715
Timing information for this millisecond binary system: $`P_p`$=5.757 ms, $`P_{orb}`$=5.741 days, $`\tau `$ (intrinsic characteristic age of pulsar) = 4.4 – 4.91 Gyrs, mass function f(M)= 1.239$`10^3`$$`M_{}`$ (Johnston et al. 1993; Bell et al. 1995). Hansen & Phinney (1998b) have discussed the evolutionary stage of this system using their own cooling models described in HP98. They found consistent solution for all masses in the range 0.15 – 0.375$`M_{}`$ with thick (in the terminology of HP98) hydrogen envelopes of 3$`\times 10^4M_{}`$.
Timing measurements by Sandhu et al. (1997) have detected a rate of change in the projected orbital separation $`a\mathrm{sin}i`$, which they interpret as a change in $`i`$ and they calculate for an upper limit for $`i`$$`<`$ $`43^0`$ and new lower limit to the mass of the companion of M$``$ 0.22$`M_{}`$. Our calculations also allow us to produce the orbital parameters and secondary mass for the PSR J0437–4715 system and fit its cooling age (2.5–5.3 Gyrs, Hansen & Phinney, 1998b), and we find that the secondary fills its Roche lobe when the orbital period $`P_i`$ is $``$ 2.5 days (Tables 1, 4). From our cooling tracks for a binary orbital period of 5.741 days, the mass of the companion is 0.21$`\pm 0.01M_{}`$ and its cooling age 1.26–2.25 Gyrs (for a Population I chemical composition). These cooling models usually have one strong (with RLOF) hydrogen shell flash, after which the helium WD enters the normal cooling phase.
### 8.2 PSR J1012+5307
Lorimer et al. (1995) determined a characteristic age of the radio pulsar to be 7 Gyr, which could be even larger if the pulsar has a significant transverse velocity (Hansen & Phinney 1998b). Using the IT86 cooling sequences, they estimated the companion to be at most 0.3 Gyr old. HP98 models yield the following results for this system: the companion mass lies in the range 0.13–0.21$`M_{}`$ and the WD age is $`<`$ 0.6 Gyr, the neutron star mass in the range 1.3–2.1$`M_{}`$.
Alberts et al. (1996) were the first to show that the cooling timescale of a low–mass WD can be substantially larger if there are no thermal flashes which lead to RLOF and a reduction of the hydrogen envelope mass. Our and DSBH98 calculations confirmed their results that for low–mass helium WDs ($`<`$ 0.2 $`M_{}`$), indeed stationary hydrogen burning plays important role. To produce short (less that one day) orbital period systems with a low–mass helium WD and a millisecond pulsar it is necessary that the secondary fills its Roche lobe between $`P_{bif}`$ and $`P_b`$ (Ergma, Sarna & Antipova, 1998). If the initial orbital period $`P_i`$ (at RLOF) is less than $`P_{bif}`$, the binary system evolves towards short orbital periods. $`P_b`$ is another critical orbital period value. If $`P_b`$ $`<`$ $`P_i(RLOF)`$ $`<`$ $`P_{bif}`$, then a short orbital period ($`<`$ 1 day) millisecond binary pulsar with low–mass helium white dwarf may form. So the initial conditions of the formation of such systems are rather important. We calculated one extra sequence to produce a binary system with orbital parameters similar to PSR J1012+5307. Initial system: 1 + 1.4 $`M_{}`$, $`P_i`$(RLOF) = 1.35 days, Z=0.01. Final system : $`M_s`$=0.168$`M_{}`$, $`P_f`$=0.605 days, $`M_{env}`$=0.041$`M_{}`$. In Fig. 11 in the effective temperature and gravity diagram we show the cooling history of this white dwarf after detachment of the Roche lobe. The two horizontal regions are the gravity values inferred by van Kerkwijk et al. (1996) (lower) and Callanan et al. (1998) (upper). Our results are consistent with the Callanan et al. (1998) estimates. It is necessary to mention that after detachment from its Roche lobe, the outer envelope is rather helium–rich. Bergeron et al. (1991) have shown that a small amount of helium in a hydrogen–dominated envelope can mimic the effect of a larger gravity.
## 9 Discussion
The results of our evolutionary calculations differ from those of Iben & Tutukov (1986) and Driebe et al. (1998) because of the different formation scenarios for low–mass helium WDs. In IT86’s calculations a donor star fills its Roche lobe while it is on the red giant branch (i.e. has a thick convective envelope) with a well developed helium core and a thin hydrogen burning layer. They proposed that the mass transfer time scale is so short that the companion will not be able to accrete the transferred matter and will itself expand and overflow its Roche lobe. The final output is the formation of a common envelope and the result of this evolution is a close binary with a helium WD of mass $`0.298\mathrm{M}_{}`$ having a rather thin ($`1.4\times 10^3\mathrm{M}_{}`$) hydrogen–rich (X=0.5) envelope.
DSBH98 did not calculate the mass exchange phases during the red giant branch evolution in detail but they also simulated the mass–exchange episode by subjecting a red giant branch model to a sufficiently large mass loss rate. In both cases (IT86 and DSBH98) mass loss starts when the star (with a well developed helium core) is on the red giant branch.
In our calculations the Roche lobe overflow starts when the secondary has either almost exhausted hydrogen in the center of the star or has a very small helium core with a thick hydrogen burning layer. During the semi–detached evolution the mass of the helium core increases from almost nothing to final value (for more detail about evolution of such systems, see Ergma, Sarna & Antipova, 1998). This is the reason that a much thicker ($`[1.56]\times 10^2\mathrm{M}_{}`$, with X ranging from 0.30 to 0.52) hydrogen–rich envelope is left on the donor star at the moment it shrinks within the Roche lobe.
The second important point where our results differ from that of DSBH98 is that in our calculations we can produce (after the secondary detaches from its Roche lobe) final millisecond binary pulsar parameters which we compare with observational data (orbital period, spin period of ms pulsar, mass of the companion). It was shown by Joss, Rappaport & Lewis (1987) and more recently by Rappaport et al. (1995) that the evolution of a binary system initially comprising of a neutron star and a low–mass giant will end up as a wide binary containing a radio pulsar and a white dwarf in a nearly circular orbit. The relation between the white dwarf mass and orbital period (see eq. (6) in Rappaport et al. 1995) shows that if the secondary fills its Roche lobe while on the red giant branch, then for $`M_{wd}0.19M_{}`$ the final orbital period would be $``$ 5 days, which is far from observed orbital period of the binary pulsar PSR J 1012+ 5307 ($`P_{orb}`$=0.6 days).
Alberts et al. (1996), DSBH98, and the results of our calculations demonstrate clearly that especially for low–mass helium WDs ($`<`$ 0.2$`M_{}`$) stationary hydrogen burning remains an important, if not the main, energy source. HP98 and BA98 did consider nuclear burning but found it to be of little importance since their artificially chosen hydrogen envelope mass was less than some critical value, disallowing significant hydrogen burning. If we compare now the cooling curves of HP98, DSH98 with ours then there is one very important difference; they did not model the evolution of the helium WD progenitor and all their cooling models (see for example Figs. 11, 12 in HP98) start with a high $`T_{eff}`$. In our models, cooling of the helium WD starts after detachment of the secondary from its Roche lobe (DSBH98 mimic this situation with mass loss from the star). This time, the secondary (proto–white dwarf) has rather low effective temperature (see for example Fig.1). During the evolution with L approximately constant, the effective temperature increases to a maximum value, after which it decreases while still having a active hydrogen shell burning source. The evolutionary time needed for the proto–white dwarf to travel from the minimum $`T_{eff}`$ (after detachment from Roche lobe) to maximum $`T_{eff}`$ depends strongly on mass of the WD (for a smaller mass a longer evolutionary time–scale).
So for low–mass helium WDs the evolutionary prehistory plays a very important role in cooling history of the white dwarf.
## 10 Conclusion
We have performed comprehensive evolutionary calculations to produce a close binary system consisting of a NS and a low–mass helium WD.
We argue that the presence of a thick hydrogen layer changes dramatically the cooling time–scale of the helium white dwarf ($`<0.25\mathrm{M}_{}`$), compared to the previous calculations (HP98, BA98) where the mass of the hydrogen envelope was chosen as free parameter and was usually one order of magnitude less than that obtained from real binary evolution computations.
Also, we have demonstrated that using new cooling tracks we can consistently explain the evolutionary status of the binary pulsar PSR J1012+53.
Tables with cooling curves are available on http://www.camk.edu.pl/ sarna/.
## Acknowledgments
We would like to thank Dr. Katrina M. Exter for help in improving the form and text of the paper. We would like to thank our referee Dr. Peter Eggleton for very useful referee opinion. At Warsaw, this work is supported through grants 2–P03D–014–07 and 2–P03D–005–16 of the Polish National Committee for Scientific Research. Also, J.A. and E.E. acknowledges support through Estonian SF grant 2446. |
warning/0002/nucl-th0002057.html | ar5iv | text | # Radiative Corrections to Elastic Electron-Proton Scattering for Polarized Electrons
## I Introduction
The radiative correction to elastic electron-proton scattering is well-known from the work of Tsai and Mo and Tsai , and the expressions given in have been used in the analysis of almost all experimental data pertaining to medium and high energy elastic electron scattering for the past three decades. Very recently, experiments using polarized electron beams have been carried out at Jefferson Lab ; specifically, longitudinally polarized electrons were scattered from unpolarized protons ($`\stackrel{}{e}pe\stackrel{}{p}`$) and the transverse and longitudinal polarizations of the recoil protons were measured in order to obtain the ratio of the proton’s elastic electromagnetic form factors, $`G_{E_p}/G_{M_p}`$. Given that radiative corrections to elastic electron-proton scattering are generally of the order of 20% - 30% for four-momentum transfer squared in the range considered in these experiments (0.5 to 3.5 (GeV/c)<sup>2</sup>), the question arises as to whether the same radiative correction used in the case of unpolarized beams and targets can be applied in the case of polarized electron beams when the polarization of the recoil proton is measured. We show here that if the approximations inherent in the calculations developed in and given in are maintained, then the same radiative correction applies both in the case of initially polarized and unpolarized electrons. In Sec. II we present the cross section for the scattering of polarized electrons from unpolarized protons in the absence of radiative corrections. In Sec. III we give each of the matrix elements associated with the radiative correction and discuss the significant approximations that are made in to evaluate their contribution to the cross section. We then show that with these approximations the radiative corrections do not depend on the polarization of either the electron or the proton in the initial or final state.
## II Differential cross section for scattering of polarized electrons
The differential cross section for the scattering of polarized electrons from unpolarized protons can be derived using standard techniques of quantum electrodynamics. We follow the conventions of Bjorken and Drell ; the metric is defined by $`p_ip_j=ϵ_iϵ_j𝐩_i𝐩_j`$. Further, $`\alpha =e^2/4\pi =1/137.036`$; $`m`$ is the electron rest mass; $`M`$ is the target nucleus rest mass; $`\kappa `$ the anomalous magnetic moment of the proton; $`p_1`$ and $`p_3`$ the initial and final electron four-momenta, respectively; $`p_2`$ and $`p_4`$ the initial and final target nucleus four-momenta, respectively; $`q=p_1p_3=p_4p_2`$ is the four-momentum transfer to the target nucleus for elastic scattering.
For one-photon exchange, the matrix element is
$$M_0=e^2\overline{u}(p_3)\gamma ^\mu u(p_1)\frac{(i)}{q^2+iϵ}\overline{u}(p_4)\mathrm{\Gamma }_\mu (q^2)u(p_2),$$
(1)
whose magnitude squared, summed over final electron spin and averaged over initial proton spin, is
$$\left|\overline{M}_0\right|^2=\frac{1}{2}\text{Tr}\left\{\gamma ^\nu \mathrm{\Lambda }_3\gamma ^\mu \mathrm{\Lambda }_1\mathrm{\Sigma }_1\right\}\text{Tr}\left\{\mathrm{\Sigma }_4\mathrm{\Lambda }_4\mathrm{\Gamma }_\mu \mathrm{\Lambda }_2\stackrel{~}{\mathrm{\Gamma }}_\nu \right\},$$
(2)
where $`\mathrm{\Lambda }_i=(p/_i+m_i)/\left(2m_i\right)`$ and $`\mathrm{\Sigma }_i=(1+\gamma _5s/_i)/2`$ are energy and spin projection operators and
$$\mathrm{\Gamma }_\mu =F_1(q^2)\gamma _\mu +\kappa F_2(q^2)\frac{i\sigma _{\mu \alpha }q^\alpha }{2M},\text{ }(\stackrel{~}{\mathrm{\Gamma }}_\nu \gamma ^0\mathrm{\Gamma }_\nu ^{}\gamma ^0)$$
(3)
is the proton-current operator. We assume high energies for the initial and final electrons ($`ϵ_1,ϵ_3>>m`$) and large momentum transfers ($`q^2>>m^2`$). Further, we express the cross section in terms of the Sachs form factors, $`G_E(q^2)`$ and $`G_M(q^2)`$, which are defined in terms of $`F_1`$ and $`F_2`$ by
$$G_E=F_1\tau \kappa F_2,G_M=F_1+\kappa F_2$$
(4)
where $`\tau =q^2/4M^2`$. Finally, we express the spin polarization four-vectors of the initial electron and final proton, $`s_1`$ and $`s_4`$ respectively, in terms of the three-dimensional unit vectors specifying the spin direction of the particles in their respective rest frames, $`𝜻_1`$ and $`𝜻_4`$. In general, for a particle of mass $`m`$, and four-momentum $`p=(ϵ,𝐩)`$, the four-vector $`s`$ is given in terms of $`𝜻`$ by ,
$`s_0`$ $`=`$ $`{\displaystyle \frac{𝜻𝐩}{m}}`$ (5)
$`𝐬`$ $`=`$ $`𝜻+𝐩\left[{\displaystyle \frac{𝜻𝐩}{m(m+ϵ)}}\right].`$ (6)
For the initial electron we have, neglecting terms of relative order $`m/ϵ_1`$,
$$s_1hp_1/m$$
(7)
where $`h𝜻_1\widehat{𝐩}_1`$. The cross section for the scattering of high energy polarized electrons into the direction $`\theta `$ by unpolarized protons initially at rest is then
$`{\displaystyle \frac{d\sigma }{d\mathrm{\Omega }}}`$ $`=`$ $`{\displaystyle \frac{\alpha ^2ϵ_3\mathrm{cos}^2\frac{\theta }{2}}{4ϵ_1^3\mathrm{sin}^4\frac{\theta }{2}}}{\displaystyle \frac{1}{(1+\tau )}}`$ (10)
$`\times (G_E^2+\tau G_M^2+2\tau (1+\tau )G_M^2\mathrm{tan}^2{\displaystyle \frac{\theta }{2}}`$
$`+h[{\displaystyle \frac{ϵ_1+ϵ_3}{M}}\sqrt{\tau (1+\tau )}G_M^2\mathrm{tan}^2{\displaystyle \frac{\theta }{2}}𝜻_4\widehat{𝐳}2\sqrt{\tau (1+\tau )}G_MG_E\mathrm{tan}{\displaystyle \frac{\theta }{2}}𝜻_4\widehat{𝐱}])`$
where we take the unit vector $`\widehat{𝐳}`$ in the direction of $`𝐩_4`$, the unit vector $`\widehat{𝐲}`$ in the direction of $`𝐩_1\times 𝐩_3`$ (i.e., perpendicular to the scattering plane) and the unit vector $`\widehat{𝐱}`$ in the scattering plane and defined by $`\widehat{𝐱}=\widehat{𝐲}\times \widehat{𝐳}`$.
In (10), the spin-independent terms give the well-known Rosenbluth cross section. The remaining terms determine the longitudinal and perpendicular polarization of the recoil proton .
## III Radiative Corrections to Elastic Electron-Proton Scattering
In this section we consider each of the terms contributing to the radiative correction to elastic electron-proton scattering as treated in the generally used analysis given in and . We show that if one makes the approximations which are inherent to the derivation given in these references, then the radiative correction to elastic electron-proton scattering is the same for polarized and unpolarized electrons and protons.
The radiative correction is comprised of the purely elastic amplitudes (electron and proton vertex corrections, electron and proton self energies, box and crossed box diagrams, and vacuum polarization terms) and inelastic amplitudes (emission of soft bremsstrahlung photons by any of the charged particles). Let us consider each of these in turn. The cross section for emission of soft photons, $`d\sigma _{\text{brem}}`$, is simply equal to a factor which multiplies the one-photon exchange cross section, $`d\sigma `$, and that factor is independent of the spins of the electrons and protons:
$$d\sigma _{\text{brem}}=\frac{\alpha }{4\pi ^2}d\sigma ^{}\frac{d^3k}{\omega }\left(\frac{p_3}{p_3k}\frac{p_1}{p_1k}\frac{p_4}{p_4k}+\frac{p_2}{p_2k}\right)^2.$$
(11)
Consider next the radiative corrections to the purely elastic cross section. To lowest order in $`\alpha `$ these are found from the cross product of the matrix element for one-photon exchange, $`M_0`$, and the matrix elements for each of the higher order processes:
$$\left|\right|^2=\left|M_0\right|^2+2\text{Re}\left\{M_0^{}\left(M_1+M_2+\mathrm{}\right)\right\}.$$
(12)
Thus, provided the matrix elements $`M_1,M_2,`$… can be expressed as $`M_0`$ times a factor which is independent of the spin of the particles, the radiative correction for elastic scattering will factor as a spin-independent term.
The matrix element for vacuum polarization, $`M_1`$, is, after charge renormalization, related simply to the matrix element $`M_0`$ by
$$M_1=M_0\underset{i}{}\mathrm{\Pi }(q^2/m_i^2)$$
(13)
in which $`\mathrm{\Pi }(q^2/m_i^2)`$ is independent of the spins of the particles , and the sum is carried over the electron and higher mass particle-antiparticle loops.
The matrix element for the electron vertex correction, $`M_2,`$ is given by
$$M_2=e^2\overline{u}(p_3)\mathrm{\Lambda }^\mu (p_3,p_1)u(p_1)\frac{(i)}{q^2+iϵ}\overline{u}(p_4)\mathrm{\Gamma }_\mu (q^2)u(p_2)$$
(14)
where
$$\mathrm{\Lambda }^\mu (p_3,p_1)=ie^2\frac{d^4k}{(2\pi )^4}\frac{1}{k^2\lambda ^2+iϵ}\gamma ^\nu \frac{1}{(p/_3k/m+iϵ)}\gamma ^\mu \frac{1}{(p/_1k/m+iϵ)}\gamma _\nu .$$
(15)
Comparing (14) with (1) we see that if the spin-operator dependence in $`\mathrm{\Lambda }^\mu (p_3,p_1)`$ reduces to $`\gamma ^\mu `$, then $`M_2`$ will be a multiple of $`M_0`$, the factor being independent of the spins of the particles. As it stands, the integral for $`\mathrm{\Lambda }^\mu (p_3,p_1)`$ is divergent. However, if we introduce a convergence factor, $`\mathrm{\Lambda }^2/(k^2\mathrm{\Lambda }^2+iϵ)`$, in the integrand then the integration can be carried out, and, taking the limit $`\mathrm{\Lambda }\mathrm{}`$, we find that $`\mathrm{\Lambda }^\mu (p_3,p_1)`$ has the form $`G_1(q^2)\gamma ^\mu +G_2(q^2)\frac{i\sigma ^{\mu \nu }q_\nu }{2m}`$, where
$$G_1^{(e)}(q^2)=\frac{\alpha }{4\pi }\left\{2(2m^2q^2)\varphi _1(\lambda ^2)+\left(\frac{3\rho ^24m^2}{\rho \rho _1}\right)\mathrm{ln}x+\frac{1}{2}+\mathrm{ln}\left(\frac{\mathrm{\Lambda }^2}{m^2}\right)\right\}$$
(16)
and
$$G_2^{(e)}(q^2)=\frac{\alpha }{4\pi }\left\{\frac{4m^2}{\rho \rho _1}\mathrm{ln}x\right\}$$
(17)
in which
$`\varphi _1(\lambda ^2)`$ $`_{\stackrel{}{\lambda 0}}`$ $`{\displaystyle \frac{1}{\rho \rho _1}}\left\{2L\left({\displaystyle \frac{1}{x}}\right){\displaystyle \frac{\pi ^2}{6}}{\displaystyle \frac{1}{2}}\mathrm{ln}^2x+\mathrm{ln}x\mathrm{ln}\left({\displaystyle \frac{\rho ^2}{\lambda ^2}}\right)\right\},`$ (19)
$`L(z)={\displaystyle _0^z}{\displaystyle \frac{\mathrm{ln}(1t)}{t}}𝑑t,`$
with $`\rho ^2=q^2+4m^2`$, $`\rho _1^2=q^2`$, and $`x=(\rho +\rho _1)/(\rho \rho _1)=(\rho +\rho _1)^2/4m^2`$. Thus for $`q^2>>m^2`$ the term $`G_2(q^2)`$ is of order $`m^2/(q^2)`$ relative to $`G_1(q^2)`$ and hence may be neglected, so that we have $`M_2=`$ $`G_1(q^2)M_0`$. The inclusion of the self energy contribution for the electron is obtained by subtracting $`\mathrm{\Lambda }^\mu (p_1,p_1)`$ from the expression given in (15), giving
$$\stackrel{~}{M}_2=\left[G_1(q^2)G_1(0)\right]M_0$$
(20)
where, for $`q^2>>m^2,`$
$`G_1(q^2)G_1(0)`$ $`=`$ $`{\displaystyle \frac{\alpha }{2\pi }}\{{\displaystyle \frac{1}{2}}\mathrm{ln}^2\left({\displaystyle \frac{q^2}{m^2}}\right)+{\displaystyle \frac{\pi ^2}{6}}`$ (23)
$`\left[\mathrm{ln}\left({\displaystyle \frac{q^2}{m^2}}\right)1\right]\mathrm{ln}\left({\displaystyle \frac{m^2}{\lambda ^2}}\right)`$
$`+{\displaystyle \frac{3}{2}}\mathrm{ln}\left({\displaystyle \frac{q^2}{m^2}}\right)2\}.`$
Finally, we consider the proton vertex correction and the box and crossed box contributions, $`M_3`$, $`M_4,`$ and $`M_5`$, respectively. The matrix elements for these corrections are given by
$$M_3=e^2\overline{u}(p_3)\gamma ^\mu u(p_1)\frac{(i)}{q^2+iϵ}\overline{u}(p_4)\mathrm{\Lambda }_\mu (p_4,p_2)u(p_2)$$
(24)
with
$`\mathrm{\Lambda }_\mu (p_4,p_2)`$ $`=`$ $`ie^2{\displaystyle \frac{d^4k}{(2\pi )^4}\frac{1}{k^2\lambda ^2+iϵ}\mathrm{\Gamma }^\nu (k^2)\frac{1}{(p/_4k/M+iϵ)}\mathrm{\Gamma }_\mu (q^2)}`$ (26)
$`\times {\displaystyle \frac{1}{(p/_2k/M+iϵ)}}\mathrm{\Gamma }_\nu (k^2),`$
$`M_4`$ $`=`$ $`(e^2)^2{\displaystyle \frac{d^4k}{(2\pi )^4}\frac{1}{k^2\lambda ^2+iϵ}\frac{1}{(kq)^2\lambda ^2+iϵ}}`$ (29)
$`\times \left[\overline{u}(p_3)\gamma ^\nu {\displaystyle \frac{1}{p/_1k/m+iϵ}}\gamma ^\mu u(p_1)\right]`$
$`\times \left[\overline{u}(p_4)\mathrm{\Gamma }_\nu ((kq)^2){\displaystyle \frac{1}{p/_2+k/M+iϵ}}\mathrm{\Gamma }_\mu (k^2)u(p_2)\right],`$
and
$`M_5`$ $`=`$ $`(e^2)^2{\displaystyle \frac{d^4k}{(2\pi )^4}\frac{1}{k^2\lambda ^2+iϵ}\frac{1}{(kq)^2\lambda ^2+iϵ}}`$ (32)
$`\times \left[\overline{u}(p_3)\gamma ^\nu {\displaystyle \frac{1}{p/_1k/m+iϵ}}\gamma ^\mu u(p_1)\right]`$
$`\times \left[\overline{u}(p_4)\mathrm{\Gamma }_\mu (k^2){\displaystyle \frac{1}{p/_4k/M+iϵ}}\mathrm{\Gamma }_\nu ((kq)^2)u(p_2)\right].`$
In general, these matrix elements depend on the initial and final spin states, and are not proportional to $`M_0`$ times a spin independent factor.
Now consider the approximation used in to evaluate these matrix elements, which we call here the soft-photon approximation. The integrands in $`M_4`$ and $`M_5`$ have two infrared divergent factors, $`[(k^2\lambda ^2+iϵ)((kq)^2\lambda ^2+iϵ)]^1`$, and are thus peaked when either of the two exchanged photons is soft, becoming divergent when $`k0`$ or when $`kq`$. We therefore first rationalize the propagators so that all spin matrices are in the numerator and then evaluate the numerators in $`M_4`$ and $`M_5`$ at these two points (first setting $`k=0`$ and then setting $`k=q`$; note that $`\mathrm{\Gamma }_\mu (0)=\gamma _\mu `$) but make no changes to the denominators. A simple calculation using the fact that we have on-shell particles shows that in fact each of the numerators has the same value for $`k=0`$ as for $`k=q`$, viz., $`4ip_1p_2q^2M_0`$ in the case of $`M_4`$ and $`4ip_3p_2q^2M_0`$ in the case of $`M_5`$. Taking this factor outside of the integral, we are left with a scalar four-point function, independent of the particle spins:
$`M_4`$ $`=`$ $`8ie^2q^2M_0p_1p_2{\displaystyle \frac{d^4k}{(2\pi )^4}\frac{1}{k^2\lambda ^2+iϵ}\frac{1}{(kq)^2\lambda ^2+iϵ}}`$ (34)
$`\times {\displaystyle \frac{1}{(k^22kp_1+iϵ)}}{\displaystyle \frac{1}{(k^2+2kp_2+iϵ)}}`$
and
$`M_5`$ $`=`$ $`8ie^2q^2M_0p_3p_2{\displaystyle \frac{d^4k}{(2\pi )^4}\frac{1}{k^2\lambda ^2+iϵ}\frac{1}{(kq)^2\lambda ^2+iϵ}}`$ (36)
$`\times {\displaystyle \frac{1}{(k^22kp_1+iϵ)}}{\displaystyle \frac{1}{(k^22kp_4+iϵ)}}.`$
We note that in an approximation is also made in the denominators of these integrals, reducing these four-point functions to three-point functions, but this is not needed for the conclusions of the present paper.
In the case of $`M_3`$, the integrand is peaked when $`k=0`$; we therefore set $`k=0`$ in all terms of the numerator of $`M_3`$, again using the fact that we have on-shell particles, and find
$$M_3=4ie^2p_4p_2M_0\frac{d^4k}{(2\pi )^4}\frac{1}{(k^2\lambda ^2+iϵ)}\frac{1}{(k^22kp_4+iϵ)}\frac{1}{(k^22kp_2+iϵ)}.$$
(37)
With the soft-photon approximation, the proton vertex correction is a multiple of $`M_0`$, and, as with $`M_2`$, the factor is independent of the spins of the particles. Again because of the soft-photon approximation, the self-energy contribution is essentially the same as that obtained for the electron: since the virtual photon in the self-energy diagrams is assumed to be soft, its interaction with the proton is given by $`\gamma _\mu ,`$ as in the case of the electron, so that the self-energy contribution is obtained by subtracting $`\mathrm{\Lambda }_\mu (p_2,p_2)`$ from the expression given in (26).
Thus, substituting the expressions for $`M_1,`$ $`M_2,`$ $`M_3,`$ $`M_4`$ and $`M_5`$ given in (13), (20), (37), (34), (36) in (12) and adding the contribution from real soft photons (11), we see that the cross section can be written in the form:
$$d\sigma _{\text{corr}}=d\sigma \left(1+\delta \right).$$
(38)
in which the radiative correction term $`\delta `$ is independent of the spins of the particles.
###### Acknowledgements.
It is a pleasure to acknowledge stimulating conversations between M. Garçon and one of the authors (L.C.M.), which provided the impetus for this work. |
warning/0002/cond-mat0002206.html | ar5iv | text | # Experimental fitting to the bipolaronic model of the normal-state resistance of Bi2Sr2CaCu2O8 single crystalsWe would like to thank A. S. Alexandrov for a critical reading of the manuscript. This research was supported by grants from the National Sciences and Engineering Research Council (NSERC) of Canada.
## Abstract
Normal-state resistance data from Bi<sub>2</sub>Sr<sub>2</sub>CaCu<sub>2</sub>O<sub>8</sub> single crystals were fitted to the (bi)polaronic conduction model, $`R=R_0(T+\sigma _bT^2)/(1+bT)`$, with satisfactory agreement over a wide temperature range. The fluctuating conduction region is found to be much narrower than that in the usual sense, as is the case for a charged Bose-gas. We estimate the effective (bi)polaron mass to be $`4m_e`$.
We report on the normal-state resistance measurement of Bi<sub>2</sub>Sr<sub>2</sub>CaCu<sub>2</sub>O<sub>8</sub> (Bi-2212) single crystals and its fitting to the bipolaron model . We found that the data fitted to this theory over wide temperature ranges. In some samples, the fitting spans from about $`T_c+40`$ K up to 300 K. Moreover, since the fitting range extends down to low temperatures near $`T_c`$, the superconducting fluctuation region is much narrower than that in the usual sense. These results provide alternative interpretations of the peculiar behavior of the normal-state resistivity.
The samples studied in the present work were Bi-2212 single crystalline whiskers , with typical dimensions of 2 to 3 mm $`\times `$ 10 $`\times `$ 0.5 $`\mu `$m along the $`a,b,`$ and $`c`$ axes, respectively. Resistance was measured by the standard four-wire method in the $`ab`$-plane with a dc current of 0.5 $`\mu `$A.
Figure 1 shows a typical plot of resistance data. The dotted line is a linear-$`T`$ fit at high temperature. As usual, the resistance shows a downward deviation around 200 K, which is generally associated with the pseudogap phenomenon (or superconducting fluctuation). In the bipolaronic scenario, on the other hand, the normal-state dynamics is dominated by small polarons and the superconducting state is a charged (2$`e`$) Bose-Einstein condensate of bipolarons . Quantitatively, this model gives for the resistance:
$$R=R_0\frac{T+\sigma _bT^2}{1+bT},$$
(1)
where $`\sigma _b`$ is the relative boson-boson scattering cross section, and $`b`$ is related to Hall coefficient,
$$R_H=\frac{v_0}{2e(nn_L)(1+bT)},$$
(2)
with $`(nn_L)(1+bT)`$ being the number of delocalized carriers in the unit-cell volume $`v_0`$ (5.40 $`\times `$ 5.41 $`\times `$ 30.9 Å<sup>3</sup>), and $`R_0`$ is a fitting parameter.
By using the Hall coefficient data from Ref. : $`R_H(10^3`$ cm<sup>3</sup>/C) = 2.5, 2.3, 2.1, and 1.9, for $`T`$ = 150, 200, 250, and 300 K, we obtained the values for $`b`$ = 0.003 K<sup>-1</sup> and $`(nn_L)=0.076`$ from Eq. 2. This results in $`\sigma _b`$ 0.0011 to 0.0016 from initial fitting of six samples. Suppose $`\sigma _b`$ is less sensitive to oxygen doping than $`b`$, so we adopt constant $`\sigma _b=0.0015`$ in the fitting. The result is as shown by the solid line in Fig. 1. As can be seen, the data closely agrees with this model over a wide temperature range. Two features in Fig. 1 are noticeable: (1) While resistance deviation from the linear-$`T`$ fitting is evident, it is well accounted for in the polaron model; and (2) Superconducting fluctuation is seen only very close to $`T_c`$ in the polaron fitting, in contrast to the much higher onset temperature if the linear-$`T`$ criterion is used. The above analysis, together with recent experiments which showed very narrow coherent region above $`T_c`$ , does not seem to support superconducting fluctuation at high temperature ($`200`$ K). The coherence length obtained by using the Aslamazov-Larkin model agrees well with literature data. Therefore, we argue the superconducting fluctuation region in the present picture is much narrower than that in the usual sense. This point agrees with the scenario of a charged Bose-gas . A similar conclusion was also established from transport measurement on YBCO .
The fitting results for three samples were shown in Fig. 2 (solid lines), all yielded mutual agreeable $`b`$ values. The paraconducting deviation happens at $`(T_c+50)`$ K, which is otherwise much higher ($`200`$ K) as described. The numbers in the plot show the values of the fitting parameters $`R_0`$/$`R`$(300 K) and $`b`$, respectively.
To further check the validity of this approach, we calculate the effective mass of the in-plane boson mass , $`m=\pi \mathrm{}^2/wa^2`$, where the lattice constant $`a=5.4`$ Å, $`w`$ is related to $`T_c`$ through
$$T_c=\frac{2w(nn_L)}{\mathrm{log}\gamma ^2}.$$
(3)
Take $`T_c=75`$ K, $`\gamma =183`$, and $`(nn_L)=0.076`$, we have $`m=4.27\times `$ free electron mass. Thus, the result may well correspond to small polarons and to inter-site bipolarons. |
warning/0002/math0002147.html | ar5iv | text | # A complete bounded minimal cylinder in ℝ³.Research partially supported by DGICYT grant number PB97-0785
## 1 Introduction
Calabi asked if it were possible to have a complete minimal surface in $`^3`$ entirely contained in a halfspace. As a consequence of the strong halfspace theorem , no such surfaces are properly immersed. The first examples of complete orientable nonflat minimal surfaces with a bounded coordinate function were obtained by Jorge and Xavier . Their construction is based on an ingenious idea of using Runge’s theorem. Later, Brito discovered a new method to construct surfaces of this kind. Examples of complete minimal surfaces with nontrivial topology, contained in a slab of $`^3`$, were obtained by Rosenberg and Toubiana , López , Costa and Simoes and Brito , among others.
A few years ago, Nadirashvili in used Runge’s theorem in a more elaborate way to produce a complete minimal disc inside a ball in $`^3`$ (see also ).
In this paper we generalize the techniques used by Nadirashvili to obtain new examples of complete minimal surfaces inside a ball in $`^3`$, with the conformal structure of an annulus. To be more precise, we have proved the following:
###### Theorem 1
There exist an open set $`A`$ of $``$ and a complete minimal immersion, $`X:A^3`$ satisfying:
1. $`X(A)`$ is a bounded set of $`^3`$;
2. $`A`$ has the conformal type of an annulus.
This theorem is proved in Section 3.
We have obtained the immersion $`X`$ as limit of a sequence of bounded minimal annuli with boundary. To construct the sequence we require a technical lemma whose proof is exhibited in Section 4. This lemma allows us to modify the intrinsic metric of a minimal annulus around the boundary, without increasing in excess the diameter of the annulus in $`^3`$.
## 2 Background and notation
The aim of this section is to fix the principal notation used in this paper, and to summarize some results about minimal surfaces.
We denote $`D_r=\{z:|z|<r\}`$, $`S_r=\{z:|z|=r\}`$ and $`D^{}=D_1\{0\}`$. Let $`X:D^{}^3`$ be a conformal minimal immersion. Then
$$\varphi _j=\frac{X_j}{u}i\frac{X_j}{v},j=1,2,3,(z=u+iv),$$
(1)
are holomorphic functions on $`D^{}`$ with real residues at $`0`$, verifying $`_{j=1}^3\varphi _{j}^{}{}_{}{}^{2}0`$ and $`_{j=1}^3|\varphi _j|^20`$. If we define
$$f=\varphi _1i\varphi _2,g=\frac{\varphi _3}{\varphi _1i\varphi _2},$$
(2)
then $`g`$ is a meromorphic function on $`D^{}`$ that coincides with the stereographic projection of the Gauss map. The behaviour of $`f`$ is determined by the rule that $`f`$ is holomorphic on $`D^{}`$, with zeroes precisely at the poles of $`g`$, but with twice order.
Conversely, if $`f,g`$ are a holomorphic and meromorphic functions, respectively, on $`D^{}`$ such that
$$\varphi _1=\frac{f}{2}(1g^2),\varphi _2=i\frac{f}{2}(1+g^2),\varphi _3=fg,$$
(3)
are holomorphic functions on $`D^{}`$, and they have no real periods in zero, then
$$X:D^{}^3$$
$$X(z)=\mathrm{Re}_{z_0}^z(\varphi _1(w),\varphi _2(w),\varphi _3(w))𝑑w+c,z_0D^{},c^3,$$
(4)
is a conformal minimal immersion. It is usual to label $`\varphi =(\varphi _1,\varphi _2,\varphi _3)`$ as the Weierstrass representation of the immersion $`X`$. We can write the conformal metric associated to the immersion $`X`$, $`\lambda _X^2(z)<,>`$, in terms of the Weierstrass representation as follows:
$$\lambda _X(z)=\frac{1}{2}|f(z)|(1+|g(z)|^2)=\frac{\varphi (z)}{\sqrt{2}}.$$
(5)
For more details minimal surfaces we refer to .
If $`\varphi :D^{}^3`$ is holomorphic, we say that $`\varphi `$ is $`𝐳^\mathrm{𝟐}`$type if $`\varphi _j(z)=\widehat{\varphi }_j(z^2),j=1,2,3`$, where $`\widehat{\varphi }_j`$ are holomorphic functions on $`D^{}`$. When the Weierstrass representation $`\varphi `$ is a $`z^2`$type map, then $`X(z)+X(z)`$ is constant on $`D^{}`$. So, we define $`S(X)=X(z)+X(z)`$ for any one particular $`zD^{}`$.
Let $`\alpha `$ be a curve in $`D^{}`$, by $`\mathrm{length}(\alpha ,X)`$ we mean the length of $`\alpha `$ with the metric associated to immersion $`X`$. For $`TD^{}`$ we define the following distance: If $`a,bT`$ let $`\mathrm{dist}_{(X,T)}(a,b)=inf\{\mathrm{length}(\alpha ,X)|\alpha :[0,1]T,\alpha (0)=a,\alpha (1)=b\}`$. If $`AT`$, $`\mathrm{dist}_{(X,T)}(z,A)`$ means the distance between point $`z`$ and set $`A`$. Any other distance or length that we use without mentioning the metric will be associated to the Euclidean metric.
By a Polygonal Pair $`(P,Q)`$, we mean a pair of closed simple curves in $`^2`$ formed by a finite number of straight segments verifying:
* $`\overline{D_{1/3}}\mathrm{Int}(Q)\overline{\mathrm{Int}(Q)}D_{2/3}\overline{D_{2/3}}\mathrm{Int}(P)\overline{\mathrm{Int}(P)}D_1`$,
* $`zP,zP`$ and $`zQzQ`$,
where $`\mathrm{Int}(\alpha )`$ denotes the interior domain bounded by a curve $`\alpha `$, and $`\mathrm{Ext}(\alpha )`$ is the exterior domain. For a pair $`(P,Q)`$, we write $`T=\mathrm{Int}(P)\overline{\mathrm{Int}(Q)}`$. If $`\xi >0`$ is small enough, $`(P^\xi ,Q^\xi )`$ represents a new polygonal pair, parallel to $`(P,Q)`$, such that:
* the Euclidean distance in $`^2`$ from $`P`$ to $`P^\xi `$ is $`\xi `$,
* the Euclidean distance in $`^2`$ from $`Q`$ to $`Q^\xi `$ is $`\xi `$,
* the corresponding set $`T^\xi `$ associated to $`(P^\xi ,Q^\xi )`$ is contained in $`T`$.
See Figure 2 in page 2.
## 3 The proof of the theorem
In order to prove the main theorem, we need the following lemma:
###### Lemma 1
Let $`X:D^{}^3`$ be a conformal minimal immersion. Consider $`(P,Q)`$ polygonal pair, and $`\rho ,r>0`$, and $`1>k>0`$, satisfying:
1. $`(1k)\rho <\mathrm{dist}_{(X,T)}(z,S_{2/3})<\rho ,zPQ,`$
2. $`X(T)B_r=\{p^3:p<r\},`$
3. $`X(z)=\mathrm{Re}\left({\displaystyle _{2/3}^z}\varphi (w)𝑑w\right)+c,`$
where $`c^3`$ and $`\varphi :D^{}^3`$ is $`z^2`$type,
4. $`S(X)=0.`$
Then, for any $`\epsilon >0`$, and for any $`s,\xi ,k^{}>0`$ verifying:
$`(1k)\rho `$ $`<`$ $`\mathrm{dist}_{(X,T^\xi )}(z,S_{2/3})<\rho ,zP^\xi Q^\xi ,`$ (6)
$`\rho `$ $`<`$ $`(1k^{})(\rho +s),`$ (7)
$`\rho k`$ $`<`$ $`s,`$ (8)
there exist a polygonal pair $`(\stackrel{~}{P},\stackrel{~}{Q})`$ and a conformal minimal immersion $`Y:D^{}^3`$, such that:
1. $`(1k^{})(\rho +s)<\mathrm{dist}_{(Y,\stackrel{~}{T})}(z,S_{2/3})<\rho +s,z\stackrel{~}{P}\stackrel{~}{Q},`$
2. $`Y(\stackrel{~}{T})B_R,R=\sqrt{r^2+(2s)^2}+\epsilon ,`$
3. $`Y(z)=\mathrm{Re}\left({\displaystyle _{2/3}^z}\psi (w)𝑑w\right)+c^{},`$
where $`c^{}^3`$ and $`\psi :D^{}^3`$ is $`z^2`$type,
4. $`S(Y)=0,`$
5. $`YX<\epsilon `$ in $`T^\xi ,`$
6. $`T^\xi \mathrm{I}\left(\stackrel{~}{T}\right)`$ and $`\stackrel{~}{T}\mathrm{I}\left(T\right)`$, where $`\mathrm{I}\left(O\right)`$ means the topological interior of the set $`O`$.
This lemma is similar in spirit to that used by Nadirashvili in his paper. However, we have worked with non simply connected planar domains bounded by polygonal pairs. So, a period problem arises. To solve this problem we have made our Weierstrass data $`\varphi `$ a $`z^2`$-type map. Furthermore, when we take limit in the conformal structure of our minimal annuli, this structure must not degenerate. This is the reason why we have dealt with pairs of parallel annuli $`T`$ and $`T^\xi `$.
Lemma 1 is proved in Section 4.
We use the lemma to construct a sequence:
$$\chi _n=(X_n:D^{}^3,(P_n,Q_n),\epsilon _n,\xi _n,k_n),$$
where $`X_n`$ is a conformal minimal immersion, $`(P_n,Q_n)`$ is a polygonal pair, and $`\{\epsilon _n\}`$, $`\{\xi _n\}`$, $`\{k_n\}`$ are decreasing sequences of non vanishing terms that converge to zero. $`\{\chi _n\}`$ must verify:
$`(1k_n)\rho _n<\mathrm{dist}_{(X_n,T_n)}(z,S_{2/3})<\rho _n,zP_nQ_n`$, where $`\rho _n=_{i=1}^n1/i,`$
$`(1k_{n1})\rho _{n1}<\mathrm{dist}_{(X_{n1},T_{n1}^{\xi _n})}(z,S_{2/3})<\rho _{n1},zP_{n1}^{\xi _n}Q_{n1}^{\xi _n},`$
$`X_n(T_n)B_{r_n}`$, where $`r_1>1`$, and $`r_n=\sqrt{r_{n1}^2+(2/n)^2}+\epsilon _n,`$
$`S(X_n)=0,`$
$`X_n(z)=\mathrm{Re}\left({\displaystyle _{2/3}^z}\varphi ^n(w)𝑑w\right)+c_n,`$
where $`c_n^3`$ and $`\varphi ^n:D^{}^3`$ is $`z^2`$type,
$`0<k_n<1`$, $`\rho _nk_n<1/(n+1)`$, and $`\epsilon _n<1/n^2,`$
$`X_nX_{n1}<\epsilon _n`$ in $`T_{n1}^{\xi _n},`$
$`\lambda _{X_n}\alpha _n\lambda _{X_{n1}}\text{in }T_{n1}^{\xi _n},`$
where $`\{\alpha _i\}_i`$ is a sequence of real numbers<sup>1</sup><sup>1</sup>1For instance, take $`\alpha _1=\frac{1}{2}e^{1/2}`$, and $`\alpha _n=e^{1/2^n}`$, $`n>1`$. such that $`0<\alpha _i<1`$ and $`\{_{i=1}^n\alpha _i\}_n`$ converges to $`1/2`$,
$`T_n\mathrm{I}\left(T_{n1}\right),`$
$`T_{n2}^{\xi _{n1}}\mathrm{I}\left(T_{n1}^{\xi _n}\right),`$
$`T_{n1}^{\xi _n}\mathrm{I}\left(T_n\right).`$
We can take, for instance,
$$\chi _1=(X_1,(P_1,Q_1),\epsilon _1=1/2,\xi _1,k_1=1/3),$$
where $`X_1:D^{}^3`$ is given by $`X_1(u+iv)=5/2(u,v,0),`$ and $`(P_1,Q_1)`$ is a suitable polygonal pair. Suppose that we have $`\chi _1,\mathrm{},\chi _n`$.
Now, we construct the $`n+1`$ term. Choose $`k_{n+1}`$ verifying $`(𝐅_{𝐧+\mathrm{𝟏}})`$, and $`\xi _{n+1}`$ verifying $`(𝐁_{𝐧+\mathrm{𝟏}})`$ and $`(𝐉_{𝐧+\mathrm{𝟏}})`$, (the choice of $`\xi _{n+1}`$ is possible since $`\chi _n`$ satisfies $`(𝐀_𝐧)`$ and $`(𝐊_𝐧)`$). Moreover, we choose two decreasing and convergent sequences to zero, $`\{\widehat{\epsilon }_m\}`$ and $`\{\widehat{\xi }_m\}`$, with $`\widehat{\xi }_m<\xi _{n+1}`$ and $`\widehat{\epsilon }_m<1/(n+1)^2`$, $`m`$. For each $`m`$, we consider $`Y_m:D^{}^3`$ and $`(\stackrel{~}{P}_m,\stackrel{~}{Q}_m)`$ given by Lemma 1, for the data:
$$X=X_n,(P,Q)=(P_n,Q_n),k^{}=k_{n+1},k=k_n,\rho =\rho _n,r=r_n,$$
$$s=1/(n+1),\epsilon =\widehat{\epsilon }_m,\xi =\widehat{\xi }_m.$$
From Assertion 5 in the lemma, we deduce that the sequence $`\{Y_m\}`$ converges to $`X_n`$ on the space $`\text{Har}(T_n)`$ of harmonic maps from $`T_n`$ in $`^3`$. This implies that $`\{\lambda _{Y_m}\}`$ converges uniformly to $`\lambda _{X_n}`$ in $`\overline{T_n^{\xi _{n+1}}}`$, and therefore there is a $`m_0`$ such that:
$$\lambda _{Y_{m_0}}\alpha _{n+1}\lambda _{X_n}\text{in }T_n^{\xi _{n+1}}.$$
(9)
We define $`X_{n+1}=Y_{m_0}`$, $`(P_{n+1},Q_{n+1})=(\stackrel{~}{P}_{m_0},\stackrel{~}{Q}_{m_0})`$, and $`\epsilon _{n+1}=\widehat{\epsilon }_{m_0}`$. Remark that $`k_{n+1},\xi _{n+1}`$ and $`\epsilon _{n+1}`$ could be chosen sufficiently small enough so that the sequences $`\{k_i\},\{\xi _i\}`$, and $`\{\epsilon _i\}`$ decrease and converge to zero. Due to the way in which we have chosen the term $`\chi _{n+1}`$ and using Lemma 1 it is easy to check that $`\chi _{n+1}`$ verifies $`(𝐀_{𝐧+\mathrm{𝟏}}),`$ $`(𝐁_{𝐧+\mathrm{𝟏}}),`$ $`\mathrm{},`$ $`(𝐊_{𝐧+\mathrm{𝟏}})`$. This concludes the construction of the sequence $`\{\chi _i\}`$.
Now, we define
$$A=\mathrm{I}\left(\underset{n}{}T_n\right).$$
The open set $`A`$ has the following properties:
1. $`A=_nT_n^{\xi _{n+1}}`$. To prove this, first observe that Properties $`(𝐈_𝐧),(𝐉_𝐧)`$, and $`(𝐊_𝐧)`$ imply $`_nT_n^{\xi _{n+1}}A`$. On the other hand, suppose that $`zA_nT_n^{\xi _{n+1}}`$. Then $`zT_nT_n^{\xi _{n+1}}`$, $`n`$. This implies that $`zA`$, which is absurd (recall that $`A`$ is open). This contradiction proves the equality.
2. $`A`$ is an open arc-connected set.
3. $`A`$ has two connected components, one of them contains zero and the other one is not bounded. Indeed, any point of $`\overline{A}`$ could be connected with $`0`$ or $`\mathrm{}`$ by a continuous curve in $`T_n`$, if $`n`$ is large enough. Then, $`A`$ has two connected components because $`\overline{A}`$ has two arc-connected components.
Therefore, $`A`$ is a domain in $``$ such that $`\{\mathrm{}\}A`$ consists of two connected components; then $`A`$ is biholomorphic to $`\{0\},D\{0\},`$ or $`C_\vartheta =\{z:\vartheta <|z|<1\}`$ (see \[5, Theorem IV.6.9\]). But $`A`$ is a hyperbolic domain, then $`A\{0\}`$. Furthermore, $`A`$ is a subset of the annulus $`C_{1/3}`$ and a generator of the homology of $`A`$ also generates the homology of $`C_{1/3}`$. So, $`AC_\vartheta `$ for a $`\vartheta ]0,1[`$.
Let $`K`$ be a compact set, subset of $`A`$. There is a $`n_0`$ such that $`KT_{n1}^{\xi _n},n>n_0`$. From $`(𝐆_𝐧)`$, we have:
$$X_NX_{n1}<\underset{i=n}{\overset{\mathrm{}}{}}\epsilon _i<\underset{i=n}{\overset{\mathrm{}}{}}1/i^2\text{in }K,N>n>n_0.$$
Thus, the sequence of minimal immersion $`\{X_n\}`$ is a Cauchy sequence in $`\text{Har}(A)`$. So, Harnack’s theorem implies that $`\{X_n\}`$ converges in $`\text{Har}(A)`$.
Let $`X:A^3`$ be the limit of $`\{X_n\}`$. $`X`$ has the following properties:
* $`X`$ is minimal and conformal.
* $`X`$ is an immersion. Indeed, for any $`zA`$ there exists $`n`$ such that $`zT_n^{\xi _{n+1}}`$. From Property $`(𝐇_𝐢)`$, we get:
$$\lambda _{X_k}(z)\alpha _k\lambda _{X_{k1}}(z)\mathrm{}\alpha _k\mathrm{}\alpha _{n+1}\lambda _{X_n}(z)\alpha _k\mathrm{}\alpha _1\lambda _{X_n(z)},k>n.$$
Taking limit as $`k\mathrm{}`$, we deduce:
$$\lambda _X(z)\frac{1}{2}\lambda _{X_n}(z)>0,$$
(10)
and so $`X`$ is an immersion.
* $`X(A)`$ is bounded in $`^3`$. Let $`zA`$ and $`n`$ such that $`zT_n^{\xi _{n+1}}`$, then
$$X(z)X(z)X_n(z)+X_n(z)\frac{1}{2}+r_n,$$
for an $`n`$ large enough. The sequence $`\{r_n\}`$ is bounded in $``$.
* The annulus $`A`$ is complete with the metric induced by $`X`$. Indeed, if $`n`$ is large enough, and taking (10) into account, one has:
$$\mathrm{dist}_{(X,T_n^{\xi _{n+1}})}(2/3,T_n^{\xi _{n+1}})>\frac{1}{2}\mathrm{dist}_{(X_n,T_n^{\xi _{n+1}})}(2/3,T_n^{\xi _{n+1}}).$$
The right hand side of this inequality is controlled by $`(𝐁_𝐧)`$, then we infer
$$\mathrm{dist}_{(X,T_n^{\xi _{n+1}})}(2/3,T_n^{\xi _{n+1}})>\frac{1}{2}(1k_n)\rho _n.$$
The completeness is due to the fact that $`\{\frac{1}{2}(1k_n)\rho _n\}_n`$ diverges.
This completes the proof of the theorem.
## 4 Proof of the lemma
This section is devoted to proving Lemma 1. As we mentioned before, it is a generalized version of that used by Nadirashvili in and Collin and Rosenberg in . Although the proof is similar, we have introduced some new techniques which permit us to apply Nadirashvili’s methods to non simply connected planar domains.
The following proposition is a direct consequence of Runge’s theorem and plays a crucial role in this section.
###### Proposition 1
Let $`\tau >1`$ and $`E_1,E_2`$ two disjoint compact sets of $``$, such that:
* $`E_i=E_i`$, $`i=1,2`$
* $`(E_1E_2)`$ has two arc-connected component, one of them contains zero and the other one is not bounded.
Then there exists $`h:\{0\}`$, a holomorphic not null function, such that:
* $`|h1|<1/\tau `$ in $`E_1`$,
* $`|h\tau |<1/\tau `$ in $`E_2`$,
* $`h(z)=\widehat{h}(z^2)`$, where $`\widehat{h}`$ is a holomorphic function in $`\{0\}`$.
* Let $`E_i^2=\{z^2:zE_i\}`$, $`i=1,2`$. It is clear that $`E_1^2`$ and $`E_2^2`$ are disjoint, and $`(E_1^2E_2^2)`$ has two connected components, one of them contains zero and the other one is not bounded. Thanks to Runge’s theorem, for any $`ϵ>0`$ there exists a holomorphic function, $`\mu :\{0\}`$, (with pole in zero), such that:
+ $`|\mu |<ϵ`$ on $`E_1^2`$,
+ $`|\mu a|<ϵ`$ on $`E_2^2`$, where $`e^a=\tau `$.
We define $`h(z)=e^{\mu (z^2)}`$, for $`ϵ`$ small enough. Q.E.D.
The main idea in the proof of Lemma 1 is to use Proposition 1 successively over a labyrinth constructed in a neigbourhood of the boundary of $`T`$. So, we modify the intrinsic metric of our immersion near the boundary, without increasing in excess the distance in $`^3`$.
Hence, the next step is to describe some subsets of $`D^{}`$ that we use to construct the above mentioned labyrinth.
Consider $`(P,Q)`$ the polygonal pair given in the statement of Lemma 1. Let $`s`$ and $`s^{}`$ be the number of sides of $`P`$ and $`Q`$, respectively, and consider $`N`$ a non trivial multiple of $`s`$ and $`s^{}`$.
###### Remark 1
Along the proof of the lemma, a set of real positive constants, $`\{r_i,i=1,\mathrm{},13\}`$, depending on $`X,(P,Q)`$, $`k`$, $`\rho `$, $`r`$,$`\epsilon `$, $`s`$, $`\xi `$ and $`k^{}`$, will appear. It is important to note that the choice of these constants does not depend on the integer $`N`$.
Let $`r_1`$ and $`r_2`$ be a lower and an upper bound, respectively, for the length of the sides of polygons $`P^\zeta `$ and $`Q^\zeta `$, $`\zeta 2/N`$. Let $`v_1,\mathrm{},v_{2N}`$ be points in the polygon $`P`$ such that they divide each side of $`P`$ into $`\frac{2N}{s}`$ equal parts. We can transfer this partition to the polygon $`P^{2/N}`$: $`v_1^{},\mathrm{},v_{2N}^{}`$; (See Figure 2). We define the following sets:
* $`L_i=`$ the segment that joins $`v_i`$ and $`v_i^{}`$, $`i=1,\mathrm{}2N`$,
* $`P_i=P^{i/N^3},i=0,\mathrm{}2N^2`$,
* $`𝒜=_{i=0}^{N^21}\overline{\mathrm{Int}(P_{2i})\mathrm{Int}(P_{2i+1})}`$, $`\stackrel{~}{𝒜}=_{i=1}^{N^2}\overline{\mathrm{Int}(P_{2i1})\mathrm{Int}(P_{2i})}`$,
* $`R=_{i=0}^{2N^2}P_i`$,
* $`=_{i=1}^NL_{2i}`$, $`\stackrel{~}{}=_{i=0}^{N1}L_{2i+1}`$,
* $`L=𝒜`$, $`\stackrel{~}{L}=\stackrel{~}{}\stackrel{~}{𝒜}`$, and $`H=RL\stackrel{~}{L}`$,
* $`\mathrm{\Omega }_N^P=\{z\mathrm{Int}(P_0)\mathrm{Int}(P_{2N^2}):\mathrm{dist}(z,H)\text{min}\{\frac{1}{4N^3},\frac{r_1}{N^2}\}\}`$,
* $`\omega _i^1`$ is the union of the segment $`L_i`$ and those connected components of $`\mathrm{\Omega }_N^P`$ which have nonempty intersection with $`L_i`$, for $`i=1,\mathrm{},N`$. Similarly, we define $`\omega _i^2`$ as the union of the segment $`L_{N+i}`$ and those connected components of $`\mathrm{\Omega }_N^P`$ which intersect $`L_{N+i}`$, for $`i=1,\mathrm{},N`$.
* $`\varpi _i^j=\{z:\mathrm{dist}(z,\omega _i^j)<\delta \}`$ where $`j=1,2`$, $`i=1,\mathrm{},N`$, and $`\delta >0`$ is chosen in such a way that the sets $`\overline{\varpi _i^j}`$, $`j=1,2`$, $`i=1,\mathrm{},N`$, are pairwise disjoint (see Figure 2),
* Finally, define $`\omega _i=\omega _i^1\omega _i^2`$ and $`\varpi _i=\varpi _i^1\varpi _i^2`$ $`i=1,\mathrm{},N`$.
As $`P`$ is symmetric, i.e. $`P=P`$, then the construction of the above sets leads us to: $`\omega _i^1=\omega _i^2`$, $`\varpi _i^1=\varpi _i^2`$.
For the polygon $`Q`$, we define, in the same way, the sets:
$$\mathrm{\Omega }_N^Q,\omega _{N+1}^j,\mathrm{},\omega _{2N}^j,\varpi _{N+1}^j,\mathrm{},\varpi _{2N}^j,j=1,2.$$
We finally define $`\mathrm{\Omega }_N=\mathrm{\Omega }_N^P\mathrm{\Omega }_N^Q`$.
The aim of the above construction is to guarantee the following, for an $`N`$ large enough,
There is a constant $`r_3`$, such that the $`\text{diam}(\varpi _i^j)r_3/N`$.
If $`\lambda ^2<,>`$ is a metric in $`D^{}`$, conformal to the Euclidean metric verifying:
$$\lambda c\text{in }T,$$
$$\lambda cN^4\text{in }\mathrm{\Omega }_N,c^+$$
and $`\alpha `$ is a curve in $`T`$ from $`S_{2/3}`$ to the boundary of $`T`$, then the length of $`\alpha `$ with this metric is greater than $`\frac{cr_1N}{2}`$. This is a consequence of the fact that each piece of $`\alpha `$, $`\alpha _i`$, $`(i=0,\mathrm{},N^21)`$, connecting $`P_{2i}`$ with $`P_{2i+2}`$, verifies the fact that either the Euclidean length of $`\alpha _i`$ is greater than $`\frac{r_1}{2N}`$, or $`\alpha _i`$ goes through a connected component of $`\mathrm{\Omega }_N`$.
Now, our purpose is to construct, for an $`N`$ large enough, a sequence of conformal minimal immersions, $`F_0=X,F_1,\mathrm{},F_{2N}`$ in $`D^{}`$ such that:
$`F_i(z)=\mathrm{Re}\left({\displaystyle _{2/3}^z}\varphi ^i(w)𝑑w\right)+c,`$
where $`c=X(2/3)`$ and $`\varphi ^i:D^{}^3`$ is $`z^2`$type,
$`\varphi ^i(z)\varphi ^{i1}(z)1/N^2,zT\varpi _i`$,
$`\varphi ^i(z)N^{7/2},z\omega _i`$,
$`\varphi ^i(z)1/\sqrt{N},z\varpi _i`$,
$`\mathrm{dist}_{𝕊^2}(G_i(z),G_{i1}(z))<\frac{1}{N\sqrt{N}},zT\varpi _i`$, where $`\mathrm{dist}_{𝕊^2}`$ is the intrinsic distance in $`𝕊^2`$, and $`G_i`$ represents the Gauss map of the immersion $`F_i`$,
there exists a set of orthogonal coordinates in $`^3`$, $`S_i=\{e_1,e_2,e_3\}`$, and a real constant $`r_4>0`$, such that:
If $`z\overline{\varpi _i}`$ and $`F_{i1}(z)1/\sqrt{N}`$ then $`(F_{i1}(z))_\mathrm{𝟏},(F_{i1}(z))_\mathrm{𝟐}<\frac{r_4}{\sqrt{N}}F_{i1}(z)`$,
$`(F_i(z))_\mathrm{𝟑}=(F_{i1}(z))_\mathrm{𝟑},z\overline{T}`$,
where $`()_𝐤`$ is the $`k^{\text{th}}`$ coordinate function with respect to $`\{e_1,e_2,e_3\}`$.
Suppose that we have $`F_0,\mathrm{},F_{j1}`$ verifying the claims $`(\mathrm{𝐏𝟏}_𝐢),`$ $`\mathrm{},`$ $`(\mathrm{𝐏𝟔}_𝐢)`$, $`i=1,\mathrm{},j1`$, then, for an $`N`$ large enough, there are positive constants $`r_5,\mathrm{},r_9`$ such that:
$`\varphi ^{j1}r_5`$ in $`T_{k=1}^{j1}\varpi _k`$.
We easily get this from $`(\mathrm{𝐏𝟐}_𝐢)`$, for $`i=1,\mathrm{},j1`$.
$`\varphi ^{j1}r_6`$ in $`T_{k=1}^{j1}\varpi _k`$.
To obtain this property, it suffices to apply $`(\mathrm{𝐏𝟐}_𝐢)`$, $`i=1,\mathrm{},j1`$, once again.
The diameter in $`^3`$ of $`F_{j1}(\varpi _j^i)`$ is less than $`r_7/N`$.
This is a consequence of $`(\mathrm{𝐋𝟏})`$, the bound of $`\text{diam}(\varpi _j^i)`$ in $`(𝐚)`$ (page (b)), and the equation (5).
The diameter in $`𝕊^2`$ of $`G_{j1}(\varpi _j^i)`$ is less than $`r_8/\sqrt{N}`$.
Indeed, from the bound of $`\text{diam}(\varpi _j^i)`$, we have a bound of diameter of $`G_0(\varpi _j^i)`$. The bound is $`sup\{(dG_0)_p:pT\}\frac{r_3}{N}`$. From successive applications of $`(\mathrm{𝐏𝟓}_𝐢)`$, we have
$$\text{diam}(G_{j1}(\varpi _j^i))<r_8/\sqrt{N}.$$
$`S(F_{j1})r_9/N`$.
This is a consequence of $`(\mathrm{𝐏𝟏}_𝐢)`$ and $`(\mathrm{𝐏𝟐}_𝐢)`$ for $`i=1,\mathrm{},j1`$.
We are going to construct $`F_j`$. We look for a set of orthogonal coordinates in $`^3`$, $`\{e_1,e_2,e_3\}`$, and a constant $`r_{10}>0`$ such that:
If $`z\varpi _j`$ and $`F_{j1}(z)1/\sqrt{N}`$, then
$$\mathrm{}(e_3,F_{j1}(z))r_{10}/\sqrt{N}\text{or}\mathrm{}(e_3,F_{j1}(z))r_{10}/\sqrt{N},$$
$`\mathrm{}(\pm e_3,G_{j1}(z))\nu /\sqrt{N}z\varpi _j,`$
where $`\mathrm{}(a,b)[0,\pi [`$ is the angle formed by $`a`$ and $`b`$ in $`^3`$, and $`\nu >1/r_6`$. We denote
$$\text{Con}(q,r)=\{x𝕊^2:\mathrm{}(x,q)r\}.$$
Let $`g_1G_{j1}(\varpi _j^1)`$ and $`g_2G_{j1}(\varpi _j^2)`$. Taking (L4) into account, the condition (D2) holds if $`e_3`$ is chosen in $`𝕊^2R`$, where
$$R=\text{Con}(g_1,\frac{r_8+\nu }{\sqrt{N}})\left[\text{Con}(g_1,\frac{r_8+\nu }{\sqrt{N}})\right]\text{Con}(g_2,\frac{r_8+\nu }{\sqrt{N}})\left[\text{Con}(g_2,\frac{r_8+\nu }{\sqrt{N}})\right].$$
The next step is to find $`e_3𝕊^2R`$ satisfying (D1) for a suitable $`r_{10}>0`$.
To do this, we define
$$F=\{p/p:pF_{j1}(\varpi _j^1)\text{ and }p\frac{1}{\sqrt{N}}\frac{r_9}{N}\}.$$
From the diameter bound of $`F_{j1}(\varpi _j^1)`$, we have that $`F\text{Con}(q,\frac{2r_7}{\sqrt{N}r_9})`$, for any $`qF`$. Consider $`r_{10}`$ such that:
$$\frac{2(r_8+\nu )}{\sqrt{N}}+\frac{2r_7}{\sqrt{N}r_9}+\frac{2r_9}{\sqrt{N}r_9}<\frac{r_{10}}{\sqrt{N}}.$$
If $`(𝕊^2R)F\mathrm{}`$, we take $`e_3(𝕊^2R)F`$. On the other hand, if $`(𝕊^2R)F=\mathrm{}`$, we take $`e_3𝕊^2R`$ such that $`\mathrm{}(e_3,q)<\frac{2(r_8+\nu )}{\sqrt{N}}`$ for some $`qF`$.
We are going to check the property (D1) in both cases.
. Take $`z\varpi _j`$ verifying $`F_{j1}(z)1/\sqrt{N}`$. If $`z\varpi _j^1`$ then an straightforward computation leads to $`\mathrm{}(e_3,F_{j1}(z))r_{10}/\sqrt{N}`$. If $`z\varpi _j^2`$, then, taking into account that $`S(F_{j1})r_9/N`$, we have $`\frac{F_{j1}(z)}{F_{j1}(z)}F`$, and $`\mathrm{}(F_{j1}(z),F_{j1}(z))\frac{2r_9}{\sqrt{N}r_9}`$. Therefore
$$\mathrm{}(e_3,F_{j1}(z))=\mathrm{}(e_3,F_{j1}(z))\mathrm{}(e_3,F_{j1}(z))+\mathrm{}(F_{j1}(z),F_{j1}(z))$$
$$\left(\frac{2(r_8+\nu )}{\sqrt{N}}+\frac{2r_7}{\sqrt{N}r_9}\right)+\frac{2r_9}{\sqrt{N}r_9}\frac{r_{10}}{\sqrt{N}}.$$
. In this case, if $`pF`$, then $`\mathrm{}(e_3,p)\mathrm{}(e_3,q)+\mathrm{}(q,p)\frac{2(r_8+\nu )}{\sqrt{N}}+\frac{2r_7}{\sqrt{N}r_9}<\frac{r_{10}}{\sqrt{N}}`$. This proves (D1) for $`z\varpi _j^1`$. If $`z\varpi _j^2`$ the proof is the same as in Case 1.
Finally, we take $`e_1,e_2`$ such that $`S_j=\{e_1,e_2,e_3\}`$ is a set of orthogonal coordinates in $`^3`$.
Let $`(f,g)`$ be the Weierstrass data of the immersion $`F_{j1}`$ in the coordinate system $`S_j`$. Let $`h`$ be the function given by Proposition 1, for $`E_1=\overline{T}\varpi _j`$, $`E_2=\omega _j`$, and $`\tau `$ large enough in order $`N`$, as we will see later. We define $`\stackrel{~}{f}=fh`$, and $`\stackrel{~}{g}=g/h`$. Now, $`\stackrel{~}{\varphi }_k^j`$, $`k=1,2,3`$, are the function defined by (3) for $`(\stackrel{~}{f},\stackrel{~}{g})`$. Then they are holomorphic and they have no periods in zero, because they are $`z^2`$type, too. Therefore, the minimal immersion $`F_j`$ is well-defined and its expression in the set of coordinates $`S_j`$ is the following:
$$F_j(z)=\mathrm{Re}\left(_{2/3}^z\stackrel{~}{\varphi }^j(w)𝑑w\right)+F_{j1}(2/3).$$
We are now going to see that $`F_j`$ verifies the properties $`(\mathrm{𝐏𝟏}_𝐣),\mathrm{},(\mathrm{𝐏𝟔}_𝐣)`$<sup>1</sup><sup>1</sup>1Note that Claims $`(\mathrm{𝐏𝟏}_𝐣),\mathrm{},(\mathrm{𝐏𝟔}_𝐣)`$ do not depend on changes of coordinates in $`^3`$.. Claim $`(\mathrm{𝐏𝟏}_𝐣)`$ easily holds. Making some calculations, we get $`(\mathrm{𝐏𝟐}_𝐣)`$, and $`(\mathrm{𝐏𝟑}_𝐣)`$, for $`\tau `$ large enough, as follows:
$$\varphi ^j\varphi ^{j1}=\frac{1}{\sqrt{2}}\left(|f(h1)|+\left|fg^2\frac{1h}{h}\right|\right)\frac{\varphi ^{j1}}{\tau 1}\frac{sup_{\overline{T}}\varphi ^{j1}}{\tau 1}\text{in }T\varpi _j,$$
and,
$$\varphi ^j=\frac{1}{\sqrt{2}}\left(|fh|+\left|\frac{fg^2}{h}\right|\right)\frac{1}{\sqrt{2}}|f||h|\frac{1}{\sqrt{2}}\underset{\overline{T}}{sup}\{|f|\}(\tau 1)\text{in }\omega _j.$$
From $`(\mathrm{𝐃𝟐})`$, we have:
$$\frac{\mathrm{sin}(\nu /\sqrt{N})}{1+\mathrm{cos}(\nu /\sqrt{N})}|g|\frac{\mathrm{sin}(\nu /\sqrt{N})}{1\mathrm{cos}(\nu /\sqrt{N})}\text{in }\varpi _j,$$
and so
$$\varphi ^j=\frac{1}{\sqrt{2}}|fg|\left(\frac{|h|}{|g|}+\frac{|g|}{|h|}\right)\frac{2}{\sqrt{2}}|fg|2\varphi ^{j1}\frac{|g|}{1+|g|^2}$$
$$r_6\mathrm{sin}(\nu /\sqrt{N})1/\sqrt{N}\text{in }\varpi _j,$$
for an $`N`$ large enough. Therefore, the property $`(\mathrm{𝐏𝟒}_𝐣)`$ is true.
Property $`(\mathrm{𝐏𝟓}_𝐣)`$ is a consequence of the following inequality:
$$2\mathrm{sin}\left(\frac{\mathrm{dist}_{𝕊^2}(G_j(z)G_{j1}(z))}{2}\right)=G_j(z)G_{j1}(z)_^3<2|\stackrel{~}{g}(z)g(z)|=$$
$$=2|g(z)||h(z)1|2\frac{sup_{\overline{T}}|g|}{\tau }zT\varpi _j.$$
Using $`(\mathrm{𝐃𝟏})`$, we get $`(\mathrm{𝐏𝟔}\mathbf{.1}_𝐣)`$, for $`r_4=r_{10}`$. And $`(\mathrm{𝐏𝟔}\mathbf{.2}_𝐣)`$ is true because, in the coordinate system $`S_j`$, we have that:
$$\varphi _3^{j1}=fg=fh\frac{g}{h}=\varphi _3^j.$$
Hence, we have constructed the immersions $`F_0,F_1,\mathrm{},F_{2N}`$ verifying Claims $`(\mathrm{𝐏𝟏}_𝐣),\mathrm{},(\mathrm{𝐏𝟔}_𝐣)`$, $`j=1,\mathrm{},2N`$. In particular, we have:
###### Proposition 2
If $`N`$ is large enough, then $`F_{2N}`$ verifies:
1. $`\rho +s<\mathrm{dist}_{(F_{2N},T)}(z,S_{2/3}),zPQ`$,
2. $`\mathrm{dist}_{(F_{2N},T^\xi )}(z,S_{2/3})<(1k^{})(\rho +s),zP^\xi Q^\xi `$,
3. there is a $`r_{11}>0`$ such that $`F_j(z)F_{j1}(z)\frac{r_{11}}{N^2}\text{ in }T\varpi _j`$,
4. $`F_{2N}X\frac{2r_{11}}{N}\text{ in }T_{j=1}^{2N}\varpi _j`$,
5. there is a polygonal pair $`(\stackrel{~}{P},\stackrel{~}{Q})`$, such that
$$(1k^{})(\rho +s)<\mathrm{dist}_{(F_{2N},\stackrel{~}{T})}(z,S_{2/3})<\rho +s,z\stackrel{~}{P}\stackrel{~}{Q},$$
6. if $`\stackrel{~}{T}`$ is the set associated to $`(\stackrel{~}{P},\stackrel{~}{Q})`$, then $`\stackrel{~}{T}\mathrm{I}\left(T\right)`$ and $`T^\xi \mathrm{I}\left(\stackrel{~}{T}\right)`$,
7. $`F_{2N}(\stackrel{~}{T})B_{R\epsilon /2},`$ where $`R=\sqrt{r^2+(2s)^2}+\epsilon `$,
where the minimal immersion $`X`$ and the constants $`\epsilon `$, $`\rho `$, $`s`$, $`r`$ and $`\xi `$ are as in Lemma 1.
* To prove Assertion (i) notice that (L2) implies:
$$\lambda _{F_{2N}}=\frac{\varphi ^{2N}}{\sqrt{2}}\frac{r_6}{\sqrt{2}}>\frac{1}{2\sqrt{N}}\text{in }T_{k=1}^{2N}\varpi _k.$$
Taking into account $`(\mathrm{𝐏𝟒}_𝐣)`$ and $`(\mathrm{𝐏𝟐}_𝐢)`$, $`i=j+1,\mathrm{},2N`$, we have
$$\lambda _{F_{2N}}\frac{\varphi ^j\varphi ^{2N}\varphi ^j}{\sqrt{2}}\frac{1}{\sqrt{2}}\left(\frac{1}{\sqrt{N}}\frac{2}{N}\right)\frac{1}{2\sqrt{N}}\text{in each }\varpi _j.$$
From $`(\mathrm{𝐏𝟑}_𝐣)`$ and $`(\mathrm{𝐏𝟐}_𝐢)`$, $`i=j+1,\mathrm{},2N`$, we obtain
$$\lambda _{F_{2N}}\frac{\varphi ^j\varphi ^{2N}\varphi ^j}{\sqrt{2}}\frac{1}{\sqrt{2}}\left(N^{7/2}\frac{2}{N}\right)\frac{1}{2\sqrt{N}}N^4\text{in each }\omega _j.$$
Using the above three inequalities and Claim (b) in page (b) we conclude the proof of the first assertion in this proposition.
To obtain Assertion (ii), consider $`zP^\xi Q^\xi `$. From (6), there is $`\alpha `$ a curve with origin $`z`$ and ending at $`z^{}S_{2/3}`$ that verifies $`\alpha T^\xi `$ and $`\mathrm{length}(\alpha ,X)<\rho `$. As $`T^\xi T_{l=1}^{2N}\varpi _l`$ (if $`N`$ is large enough), then we can apply $`(\mathrm{𝐏𝟐}_𝐣)`$, $`j=1,\mathrm{},2N`$, to obtain $`|\mathrm{length}(\alpha ,F_{2N})\mathrm{length}(\alpha ,X)|\frac{2}{\sqrt{2}N}\mathrm{length}(\alpha )`$. Bearing in mind (L2), we get $`\mathrm{length}(\alpha ,X)\frac{r_6}{\sqrt{2}}\mathrm{length}(\alpha )`$, and then
$$|\mathrm{length}(\alpha ,F_{2N})\mathrm{length}(\alpha ,X)|\frac{2}{r_6N}\rho .$$
Therefore,
$$\mathrm{length}(\alpha ,F_{2N})<\mathrm{length}(\alpha ,X)+\frac{2}{r_6N}\rho <\rho +\frac{2}{r_6N}\rho <(1k^{})(\rho +s).$$
Now we are going to prove (iii). First observe that, if $`N`$ is large enough and $`\varpi _j`$ is a set in the labyrinth $`\mathrm{\Omega }_N`$, then it is possible to find a positive constant $`r_{11}`$, only depending on $`T`$, such that: for all $`zT\varpi _j`$ there exists a curve $`\alpha _z`$ in $`T\varpi _j`$ from $`2/3`$ to $`z`$ satisfying $`\mathrm{length}(\alpha _z)<r_{11}`$. This comes from the fact that the Euclidean diameter of $`\varpi _j`$ is uniformly bounded. Using the former, we obtain
$$F_j(z)F_{j1}(z)=\mathrm{Re}_{\alpha _z}(\varphi ^j(w)\varphi ^{j1}(w))𝑑wr_{11}\frac{1}{N^2},$$
which proves Assertion (iii).
From (iii), it is not hard to deduce (iv).
Concerning (v), we are only going to construct the polygon $`\stackrel{~}{P}`$. The other polygon $`\stackrel{~}{Q}`$ can be constructed in a similar way. Let
$$𝒮=\{zD_{2/3}:(1k^{})(\rho +s)<\mathrm{dist}_{(F_{2N},T)}(z,S_{2/3})<\rho +s\}.$$
$`𝒮`$ is a not empty open subset of $`T`$. For $`\zeta >0`$ satisfying $`(1k^{})(\rho +s)<\zeta <\rho +s`$, consider
$$𝒮_\zeta =\{zD_{2/3}:\mathrm{dist}_{(F_{2N},T)}(z,S_{2/3})=\zeta \}.$$
Since $`𝒮_\zeta `$ is a compact subset of $`𝒮`$, then there are, $`B_1,\mathrm{},B_d`$, closed balls of $`^2`$ such that $`𝒮_\zeta _{i=1}^dB_i𝒮`$. Note that $`0`$ and $`\mathrm{}`$ are in disjoint arc-connected components of $`_{i=0}^dB_i`$. Then, we can construct a polygonal line $`\stackrel{~}{P}`$ in $`_{i=0}^dB_i`$ such that $`\overline{D_{2/3}}\mathrm{Int}(\stackrel{~}{P})`$. As $`\varphi ^{2N}`$ is $`z^2`$type, we have $`\lambda _{F_{2N}}(z)=\lambda _{F_{2N}}(z)`$. This means that $`\mathrm{dist}_{F_{2N}}(z,S_{2/3})=\mathrm{dist}_{F_{2N}}(z,S_{2/3})`$, $`zD^{}`$. Therefore $`\stackrel{~}{P}`$ can be chosen in such a way that $`\stackrel{~}{P}=\stackrel{~}{P}`$, because $`𝒮=𝒮`$ and $`𝒮_\zeta =𝒮_\zeta `$.
As a consequence of Assertions (i), (ii) and (v), we obtain $`\stackrel{~}{P}\mathrm{Int}(P)`$, $`\stackrel{~}{Q}\mathrm{Ext}(Q)`$, $`P^\xi \mathrm{Int}(\stackrel{~}{P})`$ and $`Q^\xi \mathrm{Ext}(\stackrel{~}{Q})`$. And so, we have that $`\stackrel{~}{T}\mathrm{I}\left(T\right)`$ and $`T^\xi \mathrm{I}\left(\stackrel{~}{T}\right)`$, which concludes (vi).
Finally, we prove Assertion (vii). Thanks to Maximum Modulus Theorem, we only need to check that
$$F_{2N}(\stackrel{~}{P}\stackrel{~}{Q})B_{R\epsilon /2}.$$
Let $`\eta \stackrel{~}{P}\stackrel{~}{Q}`$. If $`\eta T_{j=1}^{2N}\varpi _j`$, we have:
$$F_{2N}(\eta )F_{2N}(\eta )X(\eta )+X(\eta )\frac{2r_{11}}{N}+rR\epsilon /2.$$
On the other hand, if $`\eta \varpi _j`$, $`j\{1,\mathrm{},2N\}`$, the reasoning is slightly more complicated. From (v), it is possible to find a curve $`\gamma :[0,1]T`$ such that $`\gamma (0)S_{2/3},\gamma (1)=\eta `$ and $`\mathrm{length}(\gamma ,F_{2N})\rho +s`$. We define
$$\overline{t}=sup\{t[0,1]:\gamma (t)\varpi _j\},\stackrel{~}{t}=inf\{t[0,1]:\gamma (t)P^\xi \},$$
$$\overline{\eta }=\gamma (\overline{t}),\stackrel{~}{\eta }=\gamma (\stackrel{~}{t}).$$
For an $`N`$ large enough, one has $`\varpi _j\mathrm{Int}(P)\mathrm{Int}(P^\xi )`$, and so $`\stackrel{~}{t}<\overline{t}`$. Therefore, $`\gamma `$ is left divided in three disjoint pieces: $`\gamma _1`$ from $`S_{2/3}`$ to $`\stackrel{~}{\eta }`$, $`\gamma _2`$ from $`\stackrel{~}{\eta }`$ to $`\overline{\eta }`$, and $`\gamma _3`$ from $`\overline{\eta }`$ to $`\eta `$ (see Figure 3).
To continue, we need to demonstrate the existence of a constant $`r_{12}`$, that does not depend on $`N`$, such that
$$F_j(\overline{\eta })F_j(\eta )\frac{r_{12}}{N}+2s$$
(11)
Indeed,
$$F_j(\overline{\eta })F_j(\eta )F_j(\overline{\eta })F_{2N}(\overline{\eta })+F_{2N}(\overline{\eta })F_{2N}(\eta )+F_{2N}(\eta )F_j(\eta )$$
$$2\frac{2r_{11}}{N}+F_{2N}(\overline{\eta })F_{2N}(\eta )4\frac{r_{11}}{N}+\mathrm{length}(\gamma _3,F_{2N})$$
$$4\frac{r_{11}}{N}+\rho +s\mathrm{length}(\gamma _1,F_{2N})$$
(12)
Taking into account that $`\mathrm{length}(\gamma _1,F_{2N})\rho +s`$, we reason as in Assertion (ii) and obtain
$$|\mathrm{length}(\gamma _1,F_{2N})\mathrm{length}(\gamma _1,F_0)|\frac{2}{r_6N}(\rho +s).$$
(13)
Therefore, using (12) and (13), we have:
$$F_j(\overline{\eta })F_j(\eta )4\frac{r_{11}}{N}+\rho +s\mathrm{length}(\gamma _1,F_0)+\frac{2(\rho +s)}{r_6N}$$
by (6) in the hypotheses of Lemma 1, we get:
$$4\frac{r_{11}}{N}+\rho +s(1k)\rho +\frac{2(\rho +s)}{r_6N}.$$
Thus, (11) holds for $`r_{12}=4r_{11}+\frac{2(\rho +s)}{r_6}`$.
At this point, we distinguish two cases:
+ If $`F_{j1}(\overline{\eta })<1/\sqrt{N}`$ then:
$$F_{2N}(\eta )F_{2N}(\eta )F_j(\eta )+F_j(\eta )+F_j(\overline{\eta })+F_j(\overline{\eta })F_{j1}(\overline{\eta })+F_{j1}(\overline{\eta })$$
$$\frac{2r_{11}}{N}+\frac{r_{12}}{N}+2s+\frac{r_{11}}{N^2}+\frac{1}{\sqrt{N}}R\epsilon /2,$$
for an $`N`$ large enough.
+ If $`F_{j1}(\overline{\eta })>1/\sqrt{N}`$ then:
From $`(\mathrm{𝐏𝟔}\mathbf{.2}_𝐣)`$, we have, in the set of Cartesian coordinates given by $`S_j`$,
$$|(F_j(\eta ))_\mathrm{𝟑}|=|(F_{j1}(\eta ))_\mathrm{𝟑}||(F_{j1}(\eta ))_\mathrm{𝟑}(X(\eta ))_\mathrm{𝟑}+|(X(\eta ))_\mathrm{𝟑}|\frac{2r_{11}}{N}+r.$$
Using inequality (11), the fact that $`\overline{\eta }T\varpi _j`$, Assertion (iii), and Property $`(\mathrm{𝐏𝟔}\mathbf{.1}_𝐣)`$ one has
$$((F_j(\eta ))_\mathrm{𝟏},(F_j(\eta ))_\mathrm{𝟐})((F_j(\eta ))_\mathrm{𝟏},(F_j(\eta ))_\mathrm{𝟐})((F_j(\overline{\eta }))_\mathrm{𝟏},(F_j(\overline{\eta }))_\mathrm{𝟐})+$$
$$+((F_j(\overline{\eta }))_\mathrm{𝟏},(F_j(\overline{\eta }))_\mathrm{𝟐})((F_{j1}(\overline{\eta }))_\mathrm{𝟏},(F_{j1}(\overline{\eta }))_\mathrm{𝟐})+((F_{j1}(\overline{\eta }))_\mathrm{𝟏},(F_{j1}(\overline{\eta }))_\mathrm{𝟐})$$
$$\frac{r_{12}}{N}+2s+\frac{r_{11}}{N^2}+\frac{r_4}{\sqrt{N}}F_{j1}(\overline{\eta })\frac{r_{12}}{N}+2s+\frac{r_{11}}{N^2}+\frac{r_4}{\sqrt{N}}\left(\frac{2r_{11}}{N}+r\right)2s+\frac{r_{13}}{\sqrt{N}},$$
where $`r_{13}=r_{12}+r_{11}+r_4(2r_{11}+r)`$. By Pythagoras theorem,
$$F_{2N}(\eta )F_{2N}(\eta )F_j(\eta )+F_j(\eta )$$
$$\frac{2r_{11}}{N}+\sqrt{|(F_j(\eta ))_\mathrm{𝟑}|^2+((F_j(\eta ))_\mathrm{𝟏},(F_j(\eta ))_\mathrm{𝟐})^2}<\sqrt{r^2+(2s)^2}+\epsilon /2=R\epsilon /2$$
for an $`N`$ large enough.
Q.E.D.
In order to finish the proof of the lemma, we define $`Y`$ as $`Y=F_{2N}\frac{S(F_{2N})}{2}`$. It is straightforward to check that $`Y`$ verifies all the claims in Lemma 1. |
warning/0002/hep-th0002021.html | ar5iv | text | # 1 Introduction
## 1 Introduction
The idea of string theoretic description of gauge theories is an old one . Despite of the years passed on this idea, it is still activating different research works in theoretical physics . On the other hand, in the last years our understanding about string theory is changed dramatically; a stream which is usually called the “second string revolution” . The aim of this stream is formulating of a unified string theory as a fundamental theory of the known interactions. One of the most applicable tools in the above program are D$`p`$-branes . It is conjectured that D$`p`$-branes are perturbative representation of nonperturbative (strongly coupled) string theories.
It has been known for a long time that hadron-hadron scattering processes have two different behaviours depending on the amount of momentum transfer . At large momentum transfer interactions appear as interactions between the hadron constituents, partons or quarks, and some qualitative similarities to electron-hadron scattering emerge. At high energies and small momentum transfers Regge trajectories are exchanged. Regge trajectories provide a motivation for the first stringy picture of strong interaction. However, the good fitting between the Regge trajectories and the mass of strong bound-states is yet unexplained .
Deducing the apparently different observations above discussed from a unified picture is the challenge of theoretical physics and it is tempting to search for the application of the recent string theoretic progresses in this area. In this way one may find D$`p`$-branes good tools whose dynamics may be taken as a proper effective theory for the bound-states of quarks and QCD-strings (QCD electric fluxes). To use the string theory tools for QCD-strings one should replace the string theory parameters with those of QCD in a proper way. The case here is in the reverse direction of going from early days of string theory, as the theory of strong interaction, to string theory, as the theory of gravity.
To push the above idea, in two works , taking the dynamics of D0-branes as a toy model, the potential and the scattering amplitude of two D0-branes were calculated. It is found that the potential between static D0-branes is a linear potential . Also the potential between two fast decaying D0-branes, which in the extreme limit see each other instantaneously, is calculated and the general results are found in agreement with phenomenology . The scattering amplitude of two D0-branes was calculated in based on the results of , it is shown that the cross section shows the Regge pole-expansion. Regge behaviour has been used some years ago to fit the hadron-hadron total cross section data successfully (see also for some more recent application of this behaviour).
Based on the results of and some further discussions, we argue that different aspects of Light-Front formulation of QCD may be recovered by the Matrix Quantum Mechanics of D0-branes. In this paper we consider the Matrix Quantum Mechanics resulting from dimensional reduction of $`d+1`$ dimensional pure $`U(N)`$ YM theory to $`0+1`$ dimension. A detailed procedure of constructing this matrix mechanics is presented in . In analogy with string theory ($`d=9`$ or 25), we call D0-branes the free-particles sector of the moduli space. We hope that these kinds of studies shed light on possible new relation between D-brane dynamics and gauge theories. Also we adjust our discussions to be in a reasonable contact with the known phenomenological aspects, though the exact match with experiments should not be expected at this level.
In Sec.2 we review the distinguished role of Light-Front coordinates for explaining the scaling behaviour of hadrons structure functions; the same behaviour which is taken as the consequence of point-like substructure in hadrons. In Sec.3 a short review of Matrix Quantum Mechanics of D0-branes is presented. In Sec.4 the calculation of the inter D0-branes potential will be presented. The discussion on the “whiteness” of D0-branes bound-states is given in Sec.5. In Sec.6 we deal with the problem of scattering. Sec.7 is devoted to discussions. Three issues are discussed in Sec.7: 1) large-$`N`$ limit, 2) quarks, gauge theory and gravity solutions relation and 3) non-commutativity. The discussion on the non-commutativity is on a possible justification for appearance “non-commutative” coordinates in the study of “non-Abelian” bound-states, such as bound-states of quarks and gluons.
## 2 QCD, Light-Cone And Constituent Quark Picture
Before gauge theoretic description of strong interaction, QCD, there was Constituent Quark Model (CQM) for hadrons. According to CQM a meson is just a quark-antiquark bound-state and a baryon is a three-quark one. The bound-state problem has been extensively studied for years by phenomenological inter-quark potentials to calculate various low-energy quantities. The agreement between calculated and observed quantities has been always too well to justify pursuing this approach to study hadron properties .
Presently QCD is established to be the underlying theory for strong bound-states and also it has been understood that QCD-vacuum is a very complicated medium. In low energy the coupling constant is large and so quantum fluctuations are highly excited. It means that basically “sea” of quarks and gluons have considerable contribution to the properties of hadrons. Moreover, the phenomena like confinement is believed to be direct consequences of the complex nature of the QCD-vacuum. So it seems that hadron picture of QCD is not reconcilable with any few-body picture of hadrons, like CQM (see for a good discussion on this point).
Experimentally, substructure of hadrons is probed in sufficiently large momentum transfer scatterings of a fundamental particle, e.g. an electron, in the so-called Deep Inelastic Scattering (DIS) experiments. The existence of a point-like substructure, parton or quark, is taken as the reason for “scaling” behaviour of hadron structure-functions, i.e. the absence of any “scale” is the consequence of point-like objects . Along the Bjorken’s argument, and as we recall it in below, this scaling behaviour has a simple interpretation in Light-Cone point of view on the processes which are involved in DIS. The story is the same for Feynman’s parton picture of DIS experiment and the Light-Cone Frame’s cousin, the Infinite Momentum Frame (IMF) . By this simple interpretation of scaling in Light-Cone Frame we hopefully have a constituent picture for hadrons reconcilable with QCD, and it is the reason for developing the Light-Cone formulation of QCD during the past years .
The unpolarized cross-section of DIS in the lowest order is given by <sup>1</sup><sup>1</sup>1This discussion is borrowed from and .
$`k_0^{}{\displaystyle \frac{d\sigma }{d^3k^{}}}={\displaystyle \frac{2M}{sM^2}}{\displaystyle \frac{\alpha ^2}{Q^4}}l_{\mu \nu }W^{\mu \nu },`$ (2.1)
with
$`W_{\mu \nu }(p,q)`$ $`=`$ $`{\displaystyle \frac{1}{4M}}{\displaystyle \underset{\sigma }{}}{\displaystyle \frac{d^4y}{2\pi }\mathrm{e}^{iqy}p,\sigma |[J_\mu (y),J_\nu (0)]|p,\sigma },`$ (2.2)
$`l_{\mu \nu }`$ $`=`$ $`2(k_\mu k_\nu ^{}+k_\nu k_\mu ^{}{\displaystyle \frac{1}{2}}Q^2\eta _{\mu \nu }),q=kk^{},q^2=Q^2<0.`$ (2.3)
$`M`$ and $`s`$ are the mass of the nucleon and total energy respectively. The momenta are specified in the Fig.1. Also, we define the useful parameters,
$`\nu ={\displaystyle \frac{pq}{M}},x={\displaystyle \frac{Q^2}{2M\nu }},y={\displaystyle \frac{2M\nu }{sM^2}}.`$ (2.4)
Note that parameters $`x`$ and $`y`$ are dimensionless. In the rest frame of nucleon (target) we choose the z axis to be along the virtual photon momentum then we have
$`p=(M,0,0,0),q=(\nu ,0,0,\sqrt{\nu ^2+Q^2}).`$ (2.5)
In the so-called Bjorken limit, $`Q^2\mathrm{}`$, $`\nu \mathrm{}`$ and $`x`$=fixed, we have $`q=(\nu ,0,0,\nu Mx)`$. Now the statement of Bjorken scaling is as following: Up to a kinematical coefficient, the hadronic tensor $`W_{\mu \nu }`$ depends only on the parameter $`x`$ and not on $`Q^2`$. To see this, it is convenient to use Light-Cone variables $`a^\pm =(a^0\pm a^3)/\sqrt{2}`$ with scalar product as $`ab=a^+b^{}+a^{}b^+a_Tb_T`$. Thus one writes
$`W_{\mu \nu }{\displaystyle \frac{1}{4M}}{\displaystyle 𝑑y^{}\mathrm{e}^{iq^+y^{}}𝑑y^+d^2y_T\mathrm{e}^{iq^{}y^+}p|J_\mu (y)J_\nu (0)|p}.`$ (2.6)
In the Bjorken limit we have
$`q^+Mx/\sqrt{2}=\mathrm{fixed},q^{}=(2\nu +Mx)/\sqrt{2}\sqrt{2}\nu \mathrm{}.`$ (2.7)
In this limit the integrand of (2.6) contains the rapidly oscillating factor $`\mathrm{exp}(iq^{}y^+)`$ which kills all contributions to the integral except for those where the integrand is singular. Indeed the singularity of integrand comes from the current product at $`y^+0`$. In addition due to causality the integrand vanishes for $`y^2=2y^+y^{}y_T^2<0`$. So the dominant part of the integral comes from $`y^+=y_T=0`$. It explains the Bjorken scaling, i.e., the $`q^{}`$ no longer exists at $`y^+=0`$. Now it is clear that the Light-Front coordinates play a distinguished role in the understanding of the scaling behaviour in DIS experiments. The same result is also correct for Feynman’s parton description of DIS and IMF, the experimental realization of Light-Cone Frame .
## 3 Matrix Quantum Mechanics Of D0-Branes
According to string theory, D$`p`$-branes are $`p`$ dimensional objects defined as (hyper)surfaces which can trap the ends of strings and therefore it is reasonable to take their dynamics as a proper effective theory for the bound-states of quarks and QCD-strings (QCD electric fluxes).
One of the most interesting aspects of D-brane dynamics appears in their coincident limit. In the case of coinciding $`N`$ D$`p`$-branes their dynamics are captured by a $`U(N)`$ YM theory dimensionally reduced to $`p+1`$ dimensions of D$`p`$-brane world-volume . In the case of D0-branes $`p=0`$, the above dynamics reduces to quantum mechanics of matrices, because time is the only parameter in the world-line. A detailed procedure of constructing this matrix mechanics is presented in . The bosonic Lagrangian resulted from the pure YM is <sup>2</sup><sup>2</sup>2Here we take $`d`$ arbitrary.
$`L=m_0\mathrm{Tr}({\displaystyle \frac{1}{2}}D_tX_i^2+{\displaystyle \frac{1}{4(2\pi \alpha ^{})^2}}[X_i,X_j]^2),`$
$`i,j=1,\mathrm{},d,D_t=_ti[a_0,],`$
where $`\frac{1}{2\pi \alpha ^{}}`$ and $`m_0=(l_sg_s)^1`$ are the string tension and the mass of D0-branes respectively ($`l_s=\sqrt{\alpha ^{}}`$ and $`g_s`$ are the string length and coupling, respectively). For $`N`$ D0-branes $`X`$’s are in adjoint representation of $`U(N)`$ and have the usual expansion $`X_i=x_{i(a)}T_{(a)}`$, $`(a)=1,\mathrm{},N^2`$, <sup>3</sup><sup>3</sup>3To avoid confusion we put the group indices in ( ) always..
The action (3) is invariant under the residual gauge symmetry of unreduced YM theory. The transformations are:
$`\stackrel{}{X}`$ $``$ $`\stackrel{}{X^{}}=U\stackrel{}{X}U^{},`$
$`a_0`$ $``$ $`a_0^{}=Ua_0U^{}+iU_tU^{},`$ (3.2)
where $`U`$ is an arbitrary time-dependent $`N\times N`$ unitary matrix. Under these transformations one can check that:
$`D_t\stackrel{}{X}`$ $``$ $`D_t^{}\stackrel{}{X^{}}=U(D_t\stackrel{}{X})U^{},`$
$`D_tD_t\stackrel{}{X}`$ $``$ $`D_t^{}D_t^{}\stackrel{}{X^{}}=U(D_tD_t\stackrel{}{X})U^{}.`$ (3.3)
First let us search for D0-branes in the above Lagrangian:
For each direction $`i`$ there are $`N^2`$ variables and not $`N`$ which one expects for $`N`$ particles. However there is an ansatz for the equations of motion which restricts the $`U(N)`$ basis to its $`N`$ dimensional Cartan subalgebra. This ansatz causes vanishing the potential and one finds the action of $`N`$ free particles, namely:
$`S={\displaystyle 𝑑t\underset{(a)=1}{\overset{N}{}}\frac{1}{2}m_0\dot{\stackrel{}{x}}_{(a)}^2}.`$ (3.4)
In this case the $`U(N)`$ symmetry is broken to $`U(1)^N`$ and the interpretation of $`N`$ remaining variables as the classical (relative) positions of $`N`$ particles is meaningful. The center of mass of D0-branes is represented by the trace of the $`X`$ matrices.
In the case of unbroken gauge symmetry the gauge transformations mix the entries of matrices and the interpretation of positions for D0-branes remains obscure . Even in this case the center of mass is meaningful and the ambiguity about positions only remains for the relative positions of D0-branes. In unbroken phase the $`N^2N`$ non-Cartan elements of matrices have a stringy interpretation; they govern the dynamics of low lying oscillations of strings stretched between D0-branes.
The dependences of energy eigenvalues and the size of bound-states are notable. By the scalings
$`t`$ $``$ $`g_s^{1/3}t,`$
$`a_0`$ $``$ $`g_s^{1/3}a_0,`$
$`X`$ $``$ $`g_s^{1/3}X,`$ (3.5)
one finds the relevant energy and size scales as
$`E`$ $``$ $`g_s^{1/3}/l_s,`$
$`l_{d+2}`$ $`=`$ $`g_s^{1/3}l_s.`$ (3.6)
The length scale $`l_{d+2}`$ should be the fundamental length scale of the covariant $`d+2`$ dimensional theory whose Light-Cone formulation is argued to be described by action (3) with longitudinal momentum as $`m_0`$ . So it is natural to assume in our case that $`l_{d+2}`$ (for $`d=2`$) is the inverse of the 3+1 dimensional QCD mass scale, denoted by $`\mathrm{\Lambda }_{QCD}`$ <sup>4</sup><sup>4</sup>4Due to Light-Front interpretation, our $`\mathrm{\Lambda }_{QCD}`$ differs from . There $`l_s\sqrt{\alpha ^{}}`$ is taken as $`\mathrm{\Lambda }_{QCD}^1`$.. In the weak coupling $`g_s0`$ ($`m_0l_s^1`$) one finds $`l_{d+2}l_s`$ which allows to treat the bound-states of finite number of D0-branes as point-like objects in the transverse directions of the Light-Cone Frame <sup>5</sup><sup>5</sup>5Because we admit the discrete longitudinal momentum, $`m_0`$, for finite $`N`$, we are dealing with Discrete-Light-Cone-Quantization (DLCQ) . We do not emphasize on this point later., and consequently one finds $`m_0E\frac{1}{l_{d+2}^2}`$, which shows the invariance under Lorentz transformation of this combination. As we will see in Sec.6 the masses of the intermediate states in the scattering amplitude appear as $`l_{d+2}^1`$.
## 4 Known Potentials
To calculate the effective potential between D0-branes one should find the effective action around a classical configuration. This work can be done by integrating over the quantum fluctuations in a path integral. For the diagonal classical configurations, classical representations of D0-branes, the quantum fluctuations which must be integrated over are the off-diagonal entries. This work is equivalent to integrating over the oscillations of the strings stretched between D0-branes. Because here we deal with a gauge theory, and our interest is calculation around the classical field configuration, to obtain the effective action, it is convenient to use the background field method .
To calculate the effective action we write (3) in $`d+1`$ space-time dimensions in the form (in the units $`2\pi \alpha ^{}=1`$ and after the Wick rotation $`tit`$ and $`a_0ia_0`$)
$`L`$ $`=`$ $`m_0\mathrm{Tr}\left({\displaystyle \frac{1}{4}}[X_\mu ,X_\nu ]^2\right),\mu ,\nu =0,1,\mathrm{},d,`$
$`X_0`$ $`=`$ $`i_t+a_0,S={\displaystyle L𝑑t},`$ (4.1)
where $`\mu `$ and $`\nu `$ are summed over by the Euclidean metric. The one-loop effective action of (4) has been calculated several times (e.g. see the Appendix of ) and the result can be expressed as
$`({\displaystyle 𝑑t})V(X_\mu ^{cl})={\displaystyle \frac{1}{2}}\mathrm{Tr}\mathrm{log}\left(P_\lambda ^2\delta _{\mu \nu }2iF_{\mu \nu }\right)\mathrm{Tr}\mathrm{log}\left(P_\lambda ^2\right),`$ (4.2)
with
$$P_\mu [X_\mu ^{cl},],F_{\mu \nu }[f_{\mu \nu },],f_{\mu \nu }[X_\mu ^{cl},X_\nu ^{cl}],$$
and
$`P_\lambda ^2=_t^2+{\displaystyle \underset{i=1}{\overset{d}{}}}P_i^2,`$ (4.3)
with the backgrounds $`a_0^{cl}=0`$. The second term in (4.2) is due to the ghosts associated with gauge symmetry.
### 4.1 Static Potential
Here we calculate the potential between two D0-branes at rest. The classical solution which represents two D0-branes in distance $`r`$ can be introduced as
$`X_1^{cl}={\displaystyle \frac{1}{2}}\left(\begin{array}{cc}r& 0\\ 0& r\end{array}\right),X_0^{cl}=i_t\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right),`$
$`a_0^{cl}=X_i^{cl}=0,i=2,\mathrm{},d.`$ (4.4)
So one finds
$`P_1={\displaystyle \frac{r}{2}}\mathrm{\Sigma }_3,P_0=i_t1_4,P_i=0,i=2,\mathrm{},d,`$ (4.5)
where $`\mathrm{\Sigma }_3`$ is the adjoint representation of the third Pauli matrix , $`\mathrm{\Sigma }_3=[\sigma _3,]`$. The eigenvalues of $`\mathrm{\Sigma }_3`$ are 0, 0, $`\pm 2`$.
The operator $`P_\lambda ^2`$ is found to be
$`P_\lambda ^2=_t^21_4+{\displaystyle \frac{r^2}{4}}\mathrm{\Sigma }_3^2,`$ (4.6)
which is a harmonic oscillator operator whose frequency, reintroducing $`\alpha ^{}`$, is $`\omega r/\alpha ^{}`$. The one-loop effective action can be computed <sup>6</sup><sup>6</sup>6The one-loop effective action is a good approximation for $`\omega m_0\dot{r}^2`$. It gives $`rg_sl_s\dot{r}^2`$ which for $`g_s0`$ ($`m_0l_s^1`$) is satisfied for large separations and low velocities..
$`V(r)`$ $`=`$ $`({\displaystyle \frac{d1}{2}})\mathrm{Tr}\mathrm{log}\left(P_\lambda ^2\right)`$ (4.7)
$`=`$ $`\mathrm{\hspace{0.33em}2}({\displaystyle \frac{d1}{2}}){\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{ds}{s}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑k_0\mathrm{e}^{s(k_0^2+r^2)}`$
$`+\mathrm{traces}\mathrm{independent}\mathrm{of}r,`$
where 2 is for the degeneracy in eigenvalue 4 of $`\mathrm{\Sigma }_3^2`$, and $`k_0`$ is for the eigenvalues of the operator $`i_t`$. In writing the second line we have used
$$\mathrm{ln}\left(\frac{u}{v}\right)=_0^{\mathrm{}}\frac{ds}{s}(\mathrm{e}^{sv}\mathrm{e}^{su}).$$
The integrations can be performed and one finds
$`V(r)`$ $`=`$ $`\mathrm{\hspace{0.33em}2}({\displaystyle \frac{d1}{2}}){\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{ds}{s}}({\displaystyle \frac{\pi }{s}})^{\frac{1}{2}}\mathrm{e}^{sr^2}`$ (4.8)
$`=`$ $`\mathrm{\hspace{0.33em}4}\pi ({\displaystyle \frac{d1}{2}})|r|\mathrm{}(\mathrm{independent}\mathrm{of}r).`$
The linear potential is the same of phenomenology interests, see e.g. . Also it is the same which is consistent with spin-mass Regge trajectories . By restoring the $`\alpha ^{}`$ the potential will be found to be
$`V(r)=4\pi ({\displaystyle \frac{d1}{2}}){\displaystyle \frac{|r|}{2\pi \alpha ^{}}}`$ (4.9)
which has the dimension $`length^1`$. By comparison with Regge model one can have an estimation for $`\alpha ^{}`$ . The above potential can be used to describe an effective theory for the relative dynamics of D0-branes as
$`S={\displaystyle 𝑑t\left(\frac{1}{2}\frac{m_0}{2}\dot{\stackrel{}{r}}^24\pi (\frac{d1}{2})\frac{|\stackrel{}{r}|}{2\pi \alpha ^{}}\right)},`$ (4.10)
which in the range of validity of one-loop approximation, mentioned in previous footnote, it is expected to be applicable. Also by this action one obtains the energy scale as $`E\alpha ^{2/3}m_0^{1/3}g_s^{1/3}/l_s`$, as pointed in (3). The above action describes the dynamics in Light-Cone Frame with the longitudinal momentum $`m_0`$, and recalling (3) we have $`p^+p^{}m_0Eg_s^{2/3}l_s^2l_{d+2}^2`$.
It is not hard to see that the two D0-brane interaction potential is also true for every pair inside a bound-states of D0-branes. So the effective action for $`N`$ D0-branes is found to be
$`S={\displaystyle 𝑑t\left(\frac{1}{2}m_0\underset{(a)=1}{\overset{N}{}}\dot{\stackrel{}{r}}_{(a)}^24\pi (\frac{d1}{2})\underset{(a)>(b)=1}{\overset{N}{}}\frac{|\stackrel{}{r}_{(a)}\stackrel{}{r}_{(b)}|}{2\pi \alpha ^{}}\right)}.`$ (4.11)
In a recent work , by taking the linear potential between quarks of a baryonic state in transverse directions of Light-Cone Frame, the structure functions are obtained with a good agreement with observed ones.
It is useful to relate the parameter $`1/\alpha ^{}`$ in the potential with gauge theory parameters. To do this we need a string theoretic description of gauge theory, but in the Light-Cone Frame. The nearest formulation we know for this is Light-Cone–lattice gauge theory (LClgt) . In LClgt one assumes time direction and one of the spatial directions, say $`z`$, in continuum limit. The Light-Cone variables are defined as usual $`x^\pm t\pm z`$. Other spatial directions naturally play the role of transverse directions of Light-Cone Frame, which are assumed to be on a lattice in LClgt. Due to existence of a continuous time $`x^+`$, there exists a Hamiltonian formulation of the lattice gauge theory . The relation between the linear confinement potential and gauge-lattice parameters is given by :
$`V(r){\displaystyle \frac{g_{_{YM}}^2}{a^2}}|r|,`$ (4.12)
with $`a`$ as the lattice spacing parameter in the transverse directions. Comparing this with (4.9) leads to
$`{\displaystyle \frac{1}{\alpha ^{}}}{\displaystyle \frac{g_{_{YM}}^2}{a^2}}.`$ (4.13)
### 4.2 Fast Decaying D0-Branes
<sup>7</sup><sup>7</sup>7This subsection was modified based on the crucial comment by refree of EPJ.C.
For two fast decaying D0-branes one can again calculate the above potential. This work can be done by inserting for example a Gaussian function for $`k_0`$ into the (4.7). This work is equivalent to restricting the eigenvalues of the operator $`i_t`$. Having this in mind that eigenvalues of operators ($`X,i_t`$, …) represent the information corresponding to classical values of D0-branes space-time positions <sup>8</sup><sup>8</sup>8The eigenvalues of $`i_t`$ here are different from their quantum mechanical analogue which due to the Schrodinger’s equation, are energy., we find
$`V(r)`$ $`=`$ $`2({\displaystyle \frac{d1}{2}}){\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{ds}{s}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑k_0\left({\displaystyle \frac{1}{\mathrm{\Delta }}}\mathrm{e}^{\frac{k_0^2}{\mathrm{\Delta }^2}}\right)\mathrm{e}^{s(k_0^2+r^2)}`$ (4.14)
$`=`$ $`2\sqrt{\pi }({\displaystyle \frac{d1}{2}}){\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{ds}{s}}{\displaystyle \frac{e^{sr^2}}{\sqrt{s\mathrm{\Delta }^2+1}}},`$ (4.15)
in which we assumed that the D0-branes live around time zero. The last expression is infinite, but one can show that the infinite part is $`r`$-independent. One takes:
$`{\displaystyle \frac{V(r)}{(r^2)}}=2\sqrt{\pi }({\displaystyle \frac{d1}{2}}){\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dse^{sr^2}}{\sqrt{s\mathrm{\Delta }^2+1}}},`$ (4.16)
which is finite and so infinity of $`V(r)`$ is $`r`$-independent. The last integral can not be calculated exactly, though numerical comparison with phenomenology is possible. The limit $`\mathrm{\Delta }0`$ can be calculated exactly by recalling the relation:
$$\underset{\mathrm{\Delta }0}{lim}\left(\frac{1}{\mathrm{\Delta }}e^{\frac{k_0^2}{\mathrm{\Delta }^2}}\right)=\sqrt{\pi }\delta (k_0).$$
Inserting $`\delta `$-function in (4.14) one finds:
$`V(r)2({\displaystyle \frac{d1}{2}}){\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{ds}{s}}\mathrm{e}^{s(r^2)}\mathrm{ln}r,`$ (4.17)
which the last result is after extracting the $`r`$-independent infinity. This result is already consistent with phenomenology of heavy quarks , which we know their weak decay rates grow with $`(mass)^5`$. In the extreme limit $`\mathrm{\Delta }0`$, in which the two D0-branes see each other “instantaneously”, one can take them as two D(-1)-branes (D-instantons). The dynamics of D(-1)-branes are described by the action (4) but instead of the taking $`X_0`$ as $`i_t`$ one takes $`X_0`$ as a matrix which its eigenvalues represent the “instants” which D(-1)-branes occur. So the above logarithmic result also could be obtained in D(-1)-branes calculation by taking a classical solution as
$`X_1^{cl}={\displaystyle \frac{1}{2}}\left(\begin{array}{cc}r& 0\\ 0& r\end{array}\right),X_0^{cl}=\left(\begin{array}{cc}t_0& 0\\ 0& t_0\end{array}\right),a_0^{cl}=X_i^{cl}=0,i=2,\mathrm{},d,`$ (4.18)
which represents two D(-1)-branes appeared at time $`t_0`$, in distance $`r`$.
A comment is in order: from phenomenological point of view, it is known that in some cases potentials like $`r^\xi `$, with $`\xi 0.1`$ also have produced good results . This maybe can be included to our intermediate result (4.14) or logarithmic result by recalling the numerical relation $`\mathrm{ln}rr^{\eta 0}`$, which is valid for a range of $`r`$ <sup>9</sup><sup>9</sup>9One can justify by the relation:
$$\mathrm{ln}r=\underset{\eta 0}{lim}𝑑rr^{1+\eta }=\underset{\eta 0}{lim}r^\eta /\eta $$
.
## 5 White States
To determine the color of an object its dynamics should be studied in presence of external fields. For a “white” extended object, the center of mass (c.m.) moves as a free particle in a uniform electric field. Now we want to specify the color of the D0-branes bound-states. As we will see, although our formulation for dynamics of D0-branes in external YM fields seems incomplete, but there is a reasonable statement about “whiteness” of D0-branes bound-states.
### 5.1 D0-Branes In YM Background
In classical Electrodynamics besides electromagnetic fields produced by different distributions of charges and currents, we also study the dynamics of a charged particle in regions of space where electromagnetic fields exist. There is a simple question: What are the problems arising when one studies Chromodynamics in this level?
The main problem arises when one introduces sources and matches Chromodynamics with dynamics of colored objects (for example a colored particle). In case of Electrodynamics there is a simple relation. For example the equation of motion of a charge particle with mass $`m_0`$ and charge $`q`$ is
$`m_0\ddot{\stackrel{}{x}}=q(\stackrel{}{E}+\stackrel{}{v}\times \stackrel{}{B}).`$ (5.1)
The concept of gauge invariance at this level is understood as the invariance of equations of motion under gauge transformations, i.e. field strengths are invariant under gauge transformations. Now, in the case of Chromodynamics right-hand-side is a matrix and transforms as an object in adjoint representation under gauge group transformations, as
$`\stackrel{}{E}\stackrel{}{E}^{}=U\stackrel{}{E}U^{},\stackrel{}{B}\stackrel{}{B}^{}=U\stackrel{}{B}U^{}.`$ (5.2)
So the problem arises. As it is well-known for string theorists, now we have a good candidate for non-commutative coordinates which are the coordinates of coincident D0-branes. First one may rewrite (5.1) for “matrix” coordinates as
$`m_0\ddot{\stackrel{}{X}}=q(\stackrel{}{E}+\dot{\stackrel{}{X}}\times \stackrel{}{B}),`$ (5.3)
but it is not enough to have correct behaviour for the first side under gauge transformations. Here the world-line gauge symmetry (3) of D0-brane dynamics helps us, to write the generalized Lorentz equation as <sup>10</sup><sup>10</sup>10Here we drop the commutator potential in the action of D0-branes, without any lose of generality., <sup>11</sup><sup>11</sup>11One may be easier with the symmetrized version of the magnetic part as $`\frac{1}{2}(D_t\stackrel{}{X}\times \stackrel{}{B}\stackrel{}{B}\times D_t\stackrel{}{X})`$.
$`m_0D_tD_t\stackrel{}{X}=q(\stackrel{}{E}+D_t\stackrel{}{X}\times \stackrel{}{B}).`$ (5.4)
By recalling the relation (3) one observes that both sides have the same behaviour under gauge transformations. However, it seems that the picture is not complete yet. First, it is not clear what is the Lagrangian formulation of this problem. Secondly, the precise meaning of position dependences of field strengths should be clarified (there is the same question for $`U`$, the parameter of gauge transformation).
Now, the crucial observation is the decoupling of c.m. dynamics from non-Abelian parts. It is because of trace nature of $`U(1)`$ and $`SU(N)`$ parts. As we mentioned earlier the c.m. degree of freedom is described by the $`U(1)`$ part of $`U(N)`$ . So the position and the momentum of c.m. can be obtained by a simple trace
$`\stackrel{}{x}_{c.m.}{\displaystyle \frac{1}{N}}\mathrm{Tr}\stackrel{}{X},\stackrel{}{p}_{c.m.}\mathrm{Tr}\stackrel{}{P}.`$ (5.5)
To investigate the kind and amount of the charge of an object its dynamics should be studied in absence of magnetic field ($`\stackrel{}{B}=0`$) and (for extended objects) in uniform electric field ($`\stackrel{}{E}(x)=\stackrel{}{E}_0`$). So the c.m. equation of motion is
$`m_0\ddot{\stackrel{}{x}}_{c.m.}=q\stackrel{}{E}_{(1)0},`$ (5.6)
which the subscript (1) denotes that the corresponding electric field comes from the $`U(1)`$ part of $`U(N)`$. It is understood that the dynamics of c.m. will not be affected by the non-Abelian part of gauge group. It means that the c.m. is white with respect to $`SU(N)`$. This behaviour of D0-branes bound-states is the same as that of hadrons. It means that each D0-brane feels the net effect of other D0-branes as the white-complement of its color. In other words, the field fluxes extracted from one D0-brane to the other ones are the same as the flux between a color and an anti-color, Fig.2. As we have shown in Sec.4, there is a linear potential between each two static D0-branes, which is consistent with this flux-string picture. Also, the number of D0-branes in the bound-state, $`N`$, equals to that of baryons. As we mentioned before, recently the linear potential between the constituents of baryons, in the transverse directions of Light-Cone Frame, has been used successfully to obtain the structure functions.
As the final note in this part, we remind that the dynamics presented by (5.1) can be taken as for a massless particle in transverse directions in Light-Cone Frame with longitudinal momentum, $`p^+m_0`$. The fields $`\stackrel{}{E}`$ and $`\stackrel{}{B}`$ are electromagnetic fields in transverse directions. We present the derivation of this in the Appendix.
## 6 Scattering Amplitude
As a consequence of asymptotic freedom, in a suddenly collision process quarks or partons are assumed to be free. So the probe, an electron or another quark, only interacts with the hadron constituents instead of the hadron as a whole . It is the same mechanism which results scaling behaviour in hadron structure functions.
With the above in mind it is reasonable to calculate the scattering amplitude between two individual D0-branes, to find an impression about the behaviour of the scattering amplitude of two hadrons which D0-branes are assumed as their quarks. Also it is natural to assume that this result is for high energy-elastic regime of hadron collisions.
Here we use the result of . In it is shown that the quantum travelling of D0-branes can be understood by the field theory Feynman graphs and corresponding amplitudes in the Light-Cone Frame. In the following we review the approach to calculate the amplitude.
We concentrate on the limit $`\alpha ^{}0`$. In this limit to have a finite energy one has
$`[X_i,X_j]=0,i,j,`$ (6.1)
and consequently the potential term in the action vanishes. So, D0-branes do not interact and the “classical action” reduces to the action of $`N`$ free particles. We take this classical action also in the quantum case too, it is equivalent to the assumption that two quarks in two spatially well separated hadrons do not interact with each other. Since hadrons are white one can trust this assumption. However, the above observation fails when D0-branes arrive each other. When two D0-branes come very near each other two eigenvalues of $`X_i`$ matrices will be equal and the corresponding off-diagonal elements can get non-zero values. This is the same story of gauge symmetry enhancement. The fluctuations of these off-diagonal elements are responsible for the interaction between D0-branes in bound-states.
In the coincident limit the dynamics is complicated. The relative matrix position may be taken as:
$`\stackrel{}{X}=\left(\begin{array}{cc}\stackrel{}{r}/2& \stackrel{}{Y}\\ \stackrel{}{Y}^{}& \stackrel{}{r}/2\end{array}\right),`$ (6.2)
where $`Y^{}`$ is the complex conjugate of $`Y`$. By inserting this matrix into the Lagrangian one obtains:
$`S`$ $`=`$ $`{\displaystyle }dt{\displaystyle \frac{1}{2}}((2m_0)\dot{\stackrel{}{X}}_{c.m.}^2+m_0\dot{\stackrel{}{Y}}\dot{\stackrel{}{Y}}^{}{\displaystyle \frac{m_0}{4}}{\displaystyle \frac{1}{4(2\pi \alpha ^{})^2}}(1\mathrm{cos}^2\theta )\stackrel{}{r}^^2\stackrel{}{Y}\stackrel{}{Y}^{}`$ (6.3)
$`+`$ $`{\displaystyle \frac{m_0}{2}}\dot{\stackrel{}{r}}^2+O(Y^3)),`$
with $`X_{c.m.}`$ for the center of mass and $`\theta `$ is the angle between $`\stackrel{}{r}`$ and the complex vector $`\stackrel{}{Y}`$. As it is apparent in the $`\alpha ^{}0`$ limit which is the case of our interest, the $`r`$ element do not take large values and have a small range of variation. In high-tension approximation of strings ($`\alpha ^{}0`$), one can take the separation of D0-branes a constant of order $`rg_s^{1/3}l_s`$. As is noted in Sec.3, this length is the typical size of the D0-brane bound-states. So,
$`S={\displaystyle 𝑑t\left(\frac{1}{2}(2m_0)\dot{\stackrel{}{X}}_{c.m.}^2+\frac{1}{2}m_0\dot{Y}_{}\dot{Y}_{}^{}\frac{1}{2}m_0\frac{k^2r^2}{\alpha ^2}Y_{}Y_{}^{}+\frac{1}{2}\frac{m_0}{2}\dot{\stackrel{}{r}}+\mathrm{}\right)},`$ (6.4)
where in the above $`k`$ is a numerical factor depending on $`\alpha ^{}`$ and $`g_s`$ , and $`Y_{}`$ is the part of the $`\stackrel{}{Y}`$ perpendicular to the relative distance $`\stackrel{}{r}`$. The parallel part of $`\stackrel{}{Y}`$ behaves as a free part. In $`d+1`$ dimensions of space-time the dimension of $`Y_{}`$ is $`d1`$ which shows that we are encountered with $`2\times (d1)`$ harmonic oscillators because, $`Y`$ is a complex variable. This is the same number of harmonic oscillators which appears in one-loop calculations (Sec.4). These harmonic oscillators correspond to vibrations of (oriented) open strings stretched between D0-branes. In the following we ignore the radial momentum and even the angular momentum by dropping the term $`m_0\dot{\stackrel{}{r}}^2`$ and set $`r=r_0`$ for simplicity <sup>12</sup><sup>12</sup>12Setting $`r=r_0`$ may be justified by a mean value problem in integrations over constant backgrounds in the path integral as:
$$r^{d1}𝑑rDYDY^{}\mathrm{e}^{S[r,Y,Y^{}]}DYDY^{}\mathrm{e}^{S[r_0,Y,Y^{}]}.$$
.
For two D0-branes we take the probability amplitude presented by path integral as
$`x_3,x_4;t_f|x_1,x_2;t_i={\displaystyle \mathrm{e}^S}.`$ (6.5)
Based on the previous discussion, in the $`\alpha ^{}0`$ limit for (Fig.3) graph we decompose the path-integral as the following, <sup>13</sup><sup>13</sup>13Here similar to what we have in field theory we have dropped the dis-connected graphs.,
$`x_3,x_4;t_f|x_1,x_2;t_i=\left[{\displaystyle \mathrm{e}^S}\right]_{\alpha ^{}0}={\displaystyle _{t_i}^{t_f}}𝑑T_1𝑑T_2{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑X_1𝑑X_2`$ (6.6)
$`\times \left(K_{m_0}(X_1,T_1;x_1,t_i)K_{m_0}(X_1,T_1;x_2,t_i)\right)`$
$`\times \left(K_{2m_0}(X_2,T_2;X_1,T_1)K_{oscillator}(Y_{}=0,T_2;Y_{}=0,T_1)\right)`$
$`\times \left(K_{m_0}(x_3,t_f;X_2,T_2)K_{m_0}(x_4,t_f;X_2,T_2)\right),`$
which $`K_m(y_2,t_2;y_1,t_1)`$ is the non-relativistic propagator of a free particle with mass $`m`$ between $`(y_1,t_1)`$ and $`(y_2,t_2)`$ and $`K_{oscilator}(Y_{}=0,T_2;Y_{}=0,T_1)`$ is the harmonic oscillator propagator. $`𝑑T_1𝑑T_2𝑑X_1𝑑X_2`$ is for a summation over different “Joining-Splitting” times and points. We use in $`d`$ dimensions the representations
$$K_m(y_2,t_2;y_1,t_1)=\theta (t_2t_1)\frac{1}{(2\pi )^d}d^dp\mathrm{e}^{ip(y_2y_1)\frac{ip^2(t_2t_1)}{2m}},$$
$$K_{oscilator}(Y_{}=0,T_2;Y_{}=0,T_1)=\theta (T_2T_1)\left(\frac{m_0\omega }{2\pi i\mathrm{sin}[\omega (T_2T_1)]}\right)^{d1},$$
where $`\theta (t_2t_1)`$ is the step function and $`\omega `$ is the harmonic oscillator frequency, $`\omega kr_0/\alpha ^{}kg_s^{1/3}/l_s`$. Because of complex nature of $`Y_{}`$ the power for the harmonic propagator is $`2\times \frac{d1}{2}`$.
All the above results can be translated into the momentum space ($`E_k=\frac{p_k^2}{2m_0}`$ with $`k=1,2,3,4`$):
$`p_3,p_4;t_f|p_1,p_2;t_i`$ $``$ $`\mathrm{e}^{i(E_3+E_4)t_fi(E_1+E_2)t_i}{\displaystyle \underset{a=1}{\overset{4}{}}dx_a\mathrm{e}^{i(p_1x_1+p_2x_2p_3x_3p_4x_4)}}`$ (6.7)
$`\times x_3,x_4;t_f|x_1,x_2;t_i.`$
This representation is useful to calculate the cross section. The integrals can be performed and we find
$`p_3,p_4;t_f|p_1,p_2;t_i\delta ^{(d)}(p_1+p_2p_3p_4){\displaystyle _{t_i}^{t_f}}𝑑T_1𝑑T_2\theta (T_2T_1)`$ (6.8)
$`\times \mathrm{exp}({\displaystyle \frac{i(p_1^2+p_2^2)T_1}{2m_0}})\mathrm{exp}({\displaystyle \frac{iq^2(T_2T_1)}{4m_0}})\mathrm{exp}({\displaystyle \frac{i(p_3^2+p_4^2)T_2}{2m_0}})`$
$`K_{oscillator}(Y_{}=0,T_2;Y_{}=0,T_1)`$
where $`\stackrel{}{q}=\stackrel{}{p}_1+\stackrel{}{p}_2=\stackrel{}{p}_3+\stackrel{}{p}_4`$.
To have a real scattering process let us assume
$$t_i\mathrm{},t_f\mathrm{}.$$
We put $`TT_2T_1`$ which has the range $`0T\mathrm{}`$. The integrals yield
$`p_3,p_4;\mathrm{}|p_1,p_2;\mathrm{}`$ $``$ $`\delta ^{(d)}(p_1+p_2p_3p_4)\delta ({\displaystyle \frac{p_1^2}{2m_0}}+{\displaystyle \frac{p_2^2}{2m_0}}{\displaystyle \frac{p_3^2}{2m_0}}{\displaystyle \frac{p_4^2}{2m_0}})`$ (6.9)
$`{\displaystyle _0^{\mathrm{}}}𝑑T\mathrm{e}^{\frac{iT}{4m_0}(q^22(p_1^2+p_2^2))}\left({\displaystyle \frac{m_0\omega }{\mathrm{sin}(\omega T)}}\right)^{d1}.`$
Recalling the energy-momentum relation in the Light-Cone gauge ,
$$2(p_1^2+p_2^2)\stackrel{}{q}^{\mathrm{\hspace{0.33em}2}}=2(2m_0)(\frac{p_1^2+p_2^2}{2m_0})\stackrel{}{q}^{\mathrm{\hspace{0.33em}2}}=2q^+q^{}\stackrel{}{q}^{\mathrm{\hspace{0.33em}2}}=q_\mu q^\mu q_\mu ^2,$$
we find
$`p_3,p_4,E_3,E_4;\mathrm{}|p_1,p_2,E_1,E_2;\mathrm{}`$ $``$ $`\delta ^{(d)}(p_1+p_2p_3p_4)\delta (E_1+E_2E_3E_4)`$ (6.10)
$`{\displaystyle _0^{\mathrm{}}}𝑑T\mathrm{e}^{\frac{q_\mu ^2}{4m_0}T}\left({\displaystyle \frac{m_0\omega }{\mathrm{sin}(\omega T)}}\right)^{d1}.`$
We perform a cut-off for $`T`$ in small values as $`0<ϵT\mathrm{}`$, with $`ϵ`$ be small <sup>14</sup><sup>14</sup>14This cut-off is for extracting the contribution of graphs with four-legs vertex, as $`\lambda \varphi ^4`$. From time-energy uncertainty relation, we learn that these graphs are generated by super-heavy intermediate states.. By changing the integral variables as $`\mathrm{e}^{2\omega T}=\eta `$, we have
$`p_3^\mu ,p_4^\mu ;\mathrm{}|p_1^\mu ,p_2^\mu ;\mathrm{}`$ $``$ $`\delta ^{(d)}(p_1+p_2p_3p_4)\delta (p_1^{}+p_2^{}p_3^{}p_4^{})`$ (6.11)
$`{\displaystyle \frac{(m_0\omega )^{d1}}{2\omega }}{\displaystyle _0^x}𝑑\eta \eta ^{\frac{q_\mu ^2}{8m_0\omega }+\frac{d3}{2}}(1\eta )^{d+1},`$
$``$ $`\delta ^{(d)}(p_1+p_2p_3p_4)\delta (p_1^{}+p_2^{}p_3^{}p_4^{})`$
$`{\displaystyle \frac{(m_0\omega )^{d1}}{2\omega }}B_x({\displaystyle \frac{q_\mu ^2}{8m_0\omega }}+{\displaystyle \frac{d1}{2}},d+2)`$
where $`1x=\mathrm{e}^{2\omega ϵ}`$ and $`B_x`$ is the Incomplete Beta function. The longitudinal momentum conservation trivially is satisfied. Furthermore, because of the conservation of this momentum we do not expect so-called $`t`$-channel processes here.
### 6.1 Polology
Equivalently one may use the other representation of $`K_{oscillator}`$ as
$`K_{oscillator}(Y_{}=0,T_2;Y_{}=0,T_1)={\displaystyle \underset{n}{}}0|nn|0\mathrm{e}^{iE_n(T_2T_1)},`$ (6.12)
with $`E_n`$’s as the known $`H_{oscillator}`$ eigenvalues. In this representation one finds the pole expansion :
$`p_3^\mu ,p_4^\mu ;\mathrm{}|p_1^\mu ,p_2^\mu ;\mathrm{}`$ $``$ $`\delta ^{(d)}(p_1+p_2p_3p_4)\delta (p_1^{}+p_2^{}p_3^{}p_4^{})`$ (6.13)
$`\times \underset{ϵ0^+}{lim}{\displaystyle \underset{n}{}}C_n{\displaystyle \frac{i4m_0}{q_\mu q^\mu M_n^2+iϵ}}.`$
This pole expansion also can be derived by extracting the poles of the amplitude (6.11) with the condition
$`{\displaystyle \frac{q_\mu ^2}{8m_0\omega }}+{\displaystyle \frac{d1}{2}}=n,\mathrm{with}n\mathrm{as}\mathrm{a}\mathrm{positive}\mathrm{integer}.`$ (6.14)
Hence for the mass of the intermediate bound-states we obtain
$`M_n^2={\displaystyle \frac{8k(n+\frac{d1}{2})}{(g_s^{1/3}l_s)^2}}.`$ (6.15)
We recall that the combination $`g_s^{1/3}l_s`$ is $`l_{d+2}`$, the fundamental length of $`d+2`$ dimensional theory (Sec.3 and ). The Regge pole-expansion of (6.11)-(6.15) is the phenomenological promising feature of this amplitude .
## 7 Discussion
In this section we discuss some relevant issues: 1) large-$`N`$ limit, 2) quark, gauge theory and gravity solutions relations and also 3) non-commutativity.
### 7.1 Large-$`N`$
Baryons show special properties in large-$`N`$ limit of gauge theories
* Their mass grows linearly by $`N`$.
* Their size do not depend on $`N`$. So their density goes to infinity at large-$`N`$.
* Baryon-baryon force grows proportionally with $`N`$.
These properties mainly are extracted from a Hamiltonian formulation for baryons as a bound-state of $`N`$ quarks. Based on an approximation to approach the $`N`$-body problem (Hartree approximation), the above properties can be justified for baryons at large-$`N`$.
Here we try to work out the Hamiltonian formulation, and then the above mentioned properties are followed by the same reasoning of <sup>15</sup><sup>15</sup>15Because we have considered the D0-branes in Light-Cone Frame, for $`p^+=m_0l_s^1`$, the heavy quark theory of is a good approximation for the transverse dynamics of D0-branes..
In Sec.4 the effective theory for D0-branes were obtained to be
$`S={\displaystyle 𝑑t\left(\frac{1}{2}m_0\underset{(a)=1}{\overset{N}{}}\dot{\stackrel{}{r}}_{(a)}^24\pi (\frac{d1}{2})\underset{(a)>(b)=1}{\overset{N}{}}\frac{|\stackrel{}{r}_{(a)}\stackrel{}{r}_{(b)}|}{2\pi \alpha ^{}}\right)}.`$ (7.1)
Also we have found the relation between the $`\alpha ^{}`$ parameter and the coupling constant of gauge theory by comparing it to LClgt, namely $`\frac{1}{\alpha ^{}}\frac{g_{_{YM}}^2}{a^2}`$ where $`a`$ is the lattice spacing parameter. It is known that it is more convenient to replace the coupling constant by $`\frac{g_{_{YM}}}{\sqrt{N}}`$ at large-$`N`$ . So the action in terms of new parameters is
$`S={\displaystyle 𝑑t\left(\frac{1}{2}m_0\underset{(a)=1}{\overset{N}{}}\dot{\stackrel{}{r}}_{(a)}^24\pi (\frac{d1}{2})\frac{g_{_{YM}}^2}{a^2}\frac{1}{N}\underset{(a),(b)=1}{\overset{N}{}}|\stackrel{}{r}_{(a)}\stackrel{}{r}_{(b)}|\right)},`$ (7.2)
and the associated Hamiltonian is the same used in except for the potential term, which is Coulomb one there.
Here we just check the mass of baryons at large-$`N`$. The kinetic term of c.m., $`\frac{\stackrel{}{P}^2}{Nm_0}`$, grows with $`N`$, and the net potential for each D0-brane takes a factor $`\frac{1}{2}N(N1)`$ due to pair interactions. So the potential term at large-$`N`$ grows like
$`{\displaystyle \frac{1}{2}}N(N1){\displaystyle \frac{g_{_{YM}}^2}{N}}N.`$ (7.3)
It results that the energy grows as $`EN`$ at large-$`N`$. From the point of view of Light-Cone Frame the energy is $`P^{}`$. The total longitudinal momentum of this bound-state is $`P^+=Np^+`$, where $`p^+=m_0`$ is the longitudinal momentum of one D0-brane. Consequently, the invariant mass $`M`$ is
$`M^2=2P^+P^{}\stackrel{}{P}^2N^2MN.`$ (7.4)
### 7.2 Quarks, Gauge Theory And Schwartzschild Solutions Of Gravity
D$`p`$-branes are $`p`$ dimensional Schwartzschild solutions of low energy effective field theories of string theories <sup>16</sup><sup>16</sup>16In super string theories, they are charged solutions under $`p+1`$-form field.. So any proposal for equivalence between them and quarks, or at least between their dynamics and quarks dynamics, may need justification at first. Here we recall some string theoretic related issues shortly, and also try to present (maybe) a non-string theoretic related feature then.
As mentioned, D-branes are gravity solutions. On the other hand, it is known that the dynamics of these objects are captured by a gauge theory. It is one of the closest connections between gauge theories and gravity, which has been revealed by string theory. Through this relation between the dynamics of an extended object and a gauge theory, many studies have been done to develop understanding of gauge theory dynamics. One of the recent progresses in this area is the adS/CFT correspondence , to relate gauge theory dynamics at large ’t Hooft coupling ($`\lambda =g_{_{YM}}^2N`$) to gravity in the anti-de Sitter background.
The relation between gauge theory and gravity is also studied at the level of equations of motion. Both gravity and non-Abelian gauge theories, though in different orders, have nonlinear equations of motion. It is discovered that both pure gauge theories and gauge theories with matter have Schwartzschild-like solutions . By Schwartzschild-like we mean the similarity between “connections” in gauge theories (known as gauge fields $`A_{(a)}^\mu `$) and gravity (known as connection coefficients $`\mathrm{\Gamma }_{\beta \gamma }^\alpha `$). In the case of $`SU(2)`$ gauge theory with massless scalar matter field the solution is found to be
$`A_i^{(a)}`$ $`=`$ $`ϵ_{(a)ij}{\displaystyle \frac{r^j}{g_{_{YM}}r^2}}[1K(r)],`$
$`A_0^{(a)}`$ $`=`$ $`{\displaystyle \frac{r^{(a)}}{g_{_{YM}}r^2}}J(r),`$
$`\varphi ^{(a)}`$ $`=`$ $`{\displaystyle \frac{r^{(a)}}{g_{_{YM}}r^2}}H(r),`$ (7.5)
with
$`K(r)`$ $`=`$ $`{\displaystyle \frac{Cr}{1Cr}},J(r)={\displaystyle \frac{B}{1Cr}},`$
$`H(r)`$ $`=`$ $`{\displaystyle \frac{A}{1Cr}},\mathrm{with}A^2B^2=1.`$ (7.6)
The gauge fields behaviour is comparable with connection coefficients in Schwartzschild solution as
$`\mathrm{\Gamma }_{rt}^t={\displaystyle \frac{K}{2r}}{\displaystyle \frac{1}{rK}},\mathrm{\Gamma }_{rr}^r={\displaystyle \frac{K}{2r}}{\displaystyle \frac{1}{rK}},\mathrm{with}K=2GM.`$ (7.7)
Here we just review some properties of the solution (7.2) . First, both the gauge and scalar fields are singular at the radius $`r_0=C^1`$. Further, by calculating electric and magnetic fields one sees that both are singular at $`r_0`$. Therefore a particle, which carries an $`SU(2)`$ charge, becomes confined if it crosses into the region $`r<r_0`$. The singularity of field strengths at $`r_0`$ here is different from that of gravity Schwartzschild solution, which can be removed by a suitable choice of coordinates. Based on this picture of confinement of a charge in $`r<r_0`$ region, a model for confinement of gauge theories has been presented in .
Also the solution have monopole magnetic charge. This can be seen from the generalized ’t Hooft’s field strength:
$`_{\mu \nu }=_\mu (\widehat{\varphi }^{(a)}W_\nu ^{(a)})_\nu (\widehat{\varphi }^{(a)}W_\mu ^{(a)}){\displaystyle \frac{1}{g_{_{YM}}}}ϵ^{(a)(b)(c)}\widehat{\varphi }^{(a)}(_\mu \widehat{\varphi }^{(b)})(_\nu \widehat{\varphi }^{(c)}),`$ (7.8)
with $`\widehat{\varphi }^{(a)}\varphi ^{(a)}(\varphi ^{(b)}\varphi ^{(b)})^{1/2}`$. Hence, for magnetic field we find
$`_i={\displaystyle \frac{1}{2}}ϵ_{ijk}_{ij}={\displaystyle \frac{r^i}{g_{_{YM}}r^3}},`$ (7.9)
which is the magnetic field of a point monopole with charge $`4\pi /g_{_{YM}}`$. One can also find the electric field:
$`_i=_{0i}={\displaystyle \frac{r_i}{g_{_{YM}}r}}{\displaystyle \frac{d}{dr}}{\displaystyle \frac{J(r)}{r}}={\displaystyle \frac{B(2Cr1)r_i}{g_{_{YM}}r^3(1Cr)^2}},`$ (7.10)
which at $`r\mathrm{}`$ does not have the behaviour of Prasad-Sommerfield’s solution ($`1/r^2`$); and the interpretation of a net electric charge near origin is impossible. So this solution seems more like a magnetic monopole, and its relation to a “quark” (a source of electric field) is out of reach; but here the idea of Mantonen-Olive duality, which changes the role of solitonic solutions with the fundamental objects seems considerable.
### 7.3 Why Non-Commutativity?
One of the most interesting aspects of D-branes is the non-commutativity of their relative coordinates. If the model of this paper has some relation with Nature, the question will be about a possible justification for this non-commutativity. To resolve this question one may consider the following prescription: The structure of space-time has to be in correspondence and consistent with the propagation of fields. In this way one finds the space-time coordinates as $`X_\mu `$ 4-vector which behaves like electromagnetic field $`A_\mu `$ 4-vector (spin 1) under the boost transformations. This is just the same idea of special relativity to change the concept of space-time to be consistent with the Maxwell equations.
Also in this way supersymmetry is a natural continuation of the special relativity program: Adding spin $`\frac{1}{2}`$ sector to the coordinates of space-time, as the representative of the fermions of nature. This leads one to the space-time formulation of the supersymmetric theories, and in the same way ferminos are introduced into the bosonic string theory.
Now, what may be modified if nature has non-Abelian (non-commutative) gauge fields? In the present nature non-Abelian gauge fields can not make spatially long coherent states; they are confined or too heavy. But the picture may be changed inside a hadron. In fact recent developments of string theories sound this change and it is understood that non-commutative coordinates and non-Abelian gauge fields are two sides of one coin. As we discussed, the interaction between D-branes is the result of path-integrations over fluctuations of the non-commutative parts of coordinates. It means that in this picture “non-commutative” fluctuations of space-time are the source of “non-Abelian” interactions. This picture may justify involving the non-commutative coordinates in a study of bound-states of quarks and gluons. One may summarize this idea as in the below table.
| Field | Space-Time Coordinate | Theory |
| --- | --- | --- |
| Photon $`A^\mu `$ | $`X^\mu `$ | Electrodynamics (and QED) |
| Fermion $`\psi `$ | $`\theta `$, $`\overline{\theta }`$ | Supersymmetric |
| Gluon $`A_{(a)}^\mu `$ | $`X_{(a)}^\mu `$ | Chromodymamics (and QCD) |
Acknowledgement
I am grateful to S. Parvizi for useful discussions and also comments on the manuscript, and to M. Chaichian for his comment on large-$`N`$ consideration and discussions on scattering amplitudes. M.M. Sheikh-Jabbari’s comments on the revised version are deeply acknowledged. Finally I am grateful to the referee of EPJ.C., for his/her crucial comments.
## Appendix A Particle-Electrodynamics In Light-Cone Frame
We just follow the steps of in going to Light-Cone Frame. The classical action is
$`S=m{\displaystyle _1^2}𝑑\tau \sqrt{\dot{x}^2}+q{\displaystyle A_\mu \dot{x}^\mu 𝑑\tau },`$ (A.1)
with the momentum as
$`p^\mu {\displaystyle \frac{L}{\dot{x}_\mu }}=m{\displaystyle \frac{\dot{x}^\mu }{\sqrt{\dot{x}^2}}}qA^\mu .`$ (A.2)
Consequently one finds the constraints for momenta and canonical Hamiltonian
$`(p^\mu +qA^\mu )(p_\mu +qA_\mu )=m^2,`$ (A.3)
$`H_c=p_\mu \dot{x}^\mu L0.`$ (A.4)
The total Hamiltonian will be found to be
$`H_t=\lambda ((p^\mu +qA^\mu )^2m^2),`$ (A.5)
with $`\lambda `$ as Lagrange multiplier, and canonical Poisson bracket as $`\{x^\mu ,p^\nu \}=\eta ^{\mu \nu }`$. So one finds that the dynamics has gauge symmetry (reparametrization invariance) and to find the Lagrange multiplier one should fix the gauge, by condition as $`\chi (x;\tau )0`$. Preserving gauge fixing during the time gives
$`\dot{\chi }=0={\displaystyle \frac{\chi }{\tau }}+\{\chi ,H_t\},`$ (A.6)
which gives
$`\lambda =\{\chi ,\theta \}^1{\displaystyle \frac{\chi }{\tau }},\theta (p^\mu +qA^\mu )^2m^2.`$ (A.7)
The Light-Cone gauge fixing is $`\chi =\tau x^+=0`$, and also by adding the gauge fixing for the gauge field as $`A^+=0`$ , one finds for momentum conjugate of time ($`x^+`$), i.e. Hamiltonian:
$`H=p^{}={\displaystyle \frac{(\stackrel{}{p}+q\stackrel{}{A})^2}{2p^+}}qA^{},`$ (A.8)
which here we have assumed $`m=0`$. By taking $`p^+`$ as the Newtonian mass $`m_0`$ in the transverse directions and $`A^{}`$ as $`A_0`$, one gets the Lorentz’s equation of motion (5.1) by this Hamiltonian. |
warning/0002/nlin0002008.html | ar5iv | text | # 1 Introduction
## 1 Introduction
In this paper we reexamine the classical $`A_n`$–Toda systems with the aim of showing that these lattices fall into a notable class of bihamiltonian integrable systems: those for which a distinguished set of coordinates (the so–called Darboux–Nijenhuis (DN) coordinates) allows the solution of the Hamilton–Jacobi equations associated with the Hamiltonian flows by means of an (additive) separation of variables (SoV) . In particular, we will show that such coordinates arise from the geometry of the Poisson pencil after a Hamiltonian reduction process on suitable symplectic leaves.
DN coordinates (see, e.g., ) can be naturally defined on a Poisson–Nijenhuis (PN) manifold , that is, on a $`2n`$–dimensional manifold $``$ endowed with a symplectic two–form $`\omega `$ and a $`(1,1)`$ torsion free tensor $`N`$ satisfying certain compatibility conditions. In we present and discuss an intrinsic condition to characterize those Hamiltonian systems on $``$ for which DN coordinates separate the corresponding Hamilton–Jacobi equations. Moreover, some of the connections between Hamiltonian hierarchies which satisfy a certain recursion property with respect to the tensor $`N`$, and Gel’fand–Zakharevich (GZ) systems are investigated there. This paper is devoted to frame the Toda lattices into such a scheme.
The plan is as follows: in Section 2 we sketch the main points of the abovementioned SoV theory for bihamiltonian manifold, referring to for complete proofs and a more detailed discussion. Section 3 contains a formulation of the Toda lattice within the GZ scheme, that is, taking as starting point its Poisson pencil and the problem of finding the Casimir functions. Section 4 concerns the application of the bihamiltonian SoV theory to this family of integrable systems. Finally, in Section 5 we treat the three-particle case to give a feeling of how the method works.
## 2 Separation of variables on PN manifolds
Let $`(,\omega )`$ be a $`2n`$–dimensional symplectic manifold endowed with a Nijenhuis tensor field $`N`$ compatible with $`\omega `$ (in the sense of the theory of bihamiltonian manifolds). These manifolds are called Poisson–Nijenhuis manifolds . Examples of such manifolds are provided by bihamiltonian manifolds endowed with a pair of Poisson bivectors $`(P_0,P_1)`$ one of which, say $`P_0`$, is invertible, In this case, $`\omega =P_0^1`$ and $`N=P_1P_0^1`$.
###### Definition 2.1
By Darboux–Nijenhuis coordinates on $``$ we mean any system of local coordinates $`(\lambda _j,\mu _j)_{j=1,\mathrm{},n}`$ which enjoy the following two properties:
i) $`\omega `$ takes the canonical form
$$\omega =\underset{i=1}{\overset{n}{}}d\lambda _id\mu _i;$$
ii) the adjoint Nijenhuis operator $`N^{}`$ takes the diagonal form
$$N^{}d\lambda _j=\lambda _jd\lambda _j,N^{}d\mu _j=\lambda _jd\mu _j.$$
It has been shown that DN coordinates exist on any PN manifold where $`N`$ has $`n`$ functionally independent eigenvalues. In this case the coordinates $`\lambda _j`$ can be computed algebraically as the roots of the minimal polynomial of $`N`$,
$$C(\lambda )=\text{Det}\left(N\lambda \mathrm{𝟏}\right)^{\frac{1}{2}}.$$
(2.1)
On the contrary, the complementary coordinates must be computed (in general) by a method involving quadratures.
In we characterize a class of Hamiltonians on $``$ whose associated Hamilton–Jacobi equations can be solved by separation of variables in DN coordinates. Let $`(H_1,\mathrm{},H_n)`$ be a set of functionally independent (Hamiltonian) functions that are in involution with respect to the canonical Poisson bracket defined by
$$\{f,g\}=\omega (X_f,X_g).$$
We assume that the Lagrangian foliation defined by the functions $`H_i`$ is invariant with respect to $`N`$. This is tantamount to saying that, at a generic point $`m`$, the differentials $`dH_k`$ span an $`n`$–dimensional vector subspace of $`T_m^{}`$ which is closed under the action of $`N^{}`$. Hence there exists an $`n\times n`$ matrix $`𝖥`$, whose entries are functions on $``$, such that
$$N^{}dH_i=\underset{j=1}{\overset{n}{}}𝖥_i^jdH_j.$$
(2.2)
###### Definition 2.2
The Hamiltonians $`(H_1,\mathrm{},H_n)`$ are separable in the DN coordinates if there exists an $`n\times n`$ invertible matrix $`𝖳`$ and an $`n`$–component vector $`V`$ such that
$$𝖳H=V,$$
(2.3)
where $`H=(H_1,\mathrm{},H_n)^T`$, and the matrix $`𝖳`$ and the vector $`V`$ possess the Stäckel properties:
1. the entries of the $`j^{th}`$ row of $`𝖳`$ depend only on the conjugated coordinates $`(\lambda _j,\mu _j)`$.
2. the $`j^{th}`$ component of the vector $`V`$ depends only on $`(\lambda _j,\mu _j)`$ as well.
A remarkable “separability test” is given by the following
###### Theorem 2.3
The Hamiltonians $`(H_1,\mathrm{},H_n)`$ are separable if and only if the matrix $`𝖥`$ verifies the equation
$$N^{}d𝖥=𝖥d𝖥.$$
(2.4)
Two remarks are in order to explain this theorem. First of all, equation (2.4) must be read as follows. In the left hand–side, $`d𝖥`$ is the matrix whose entries are the differentials of the entries of $`𝖥`$, and $`N^{}`$ acts separately on each entry. Secondly, one should notice that (2.4) is a coordinate free test of separability, that can be checked without computing the DN coordinates. Once the test is passed one can construct the Stäckel matrix $`𝖳`$, still in general coordinates, by a simple algebraic procedure. One has to consider the eigenvectors of the matrix $`𝖥`$ and form with them a (suitably normalized) matrix $`𝖳`$ that diagonalizes $`𝖥`$:
$$𝖥=𝖳^1\mathrm{\Lambda }𝖳,\mathrm{\Lambda }=\text{ diag }(\lambda _1,\mathrm{},\lambda _n).$$
By condition (2.4), this matrix is a Stäckel matrix; by condition (2.2), the vector $`V=𝖳H`$ verifies the Stäckel property. Then, once constructed the DN coordinates, the Hamilton–Jacobi equations associated with $`(H_1\mathrm{},H_n)`$ can be easily solved by separation of variables. Notice that the DN coordinates separate at once the HJ equations associated with any of the Hamiltonians $`H_i`$.
To complete the construction of the DN coordinates, that is, to construct algebraically the coordinates $`\mu _j`$ conjugated to the eigenvalues $`\lambda _j`$ of $`N`$, the following procedure is often useful. We consider the Hamiltonian vector field $`Y`$ associated (by the symplectic form $`\omega `$) with the function $`\frac{1}{2}\text{Tr}(N)`$, and the space of functions $`F(x;\lambda )`$, depending smoothly on $`x`$ and holomorphically on the parameter $`\lambda `$. We denote with $`F(x;\lambda _j)`$ the evaluation of $`F(x;\lambda )`$ at $`\lambda =\lambda _j`$. If $`N^{}dF(x;\lambda _j)=\lambda _jdF(x;\lambda _j)`$ for all $`j=1,\mathrm{},n`$, we say that $`F(x;\lambda )`$ is an exact eigenvector of $`N^{}`$.
###### Theorem 2.4
If $`F(x;\lambda )`$ is an exact eigenvector of $`N^{}`$, satisfying the “normalization property” $`Y(F(x;\lambda ))=1`$, then the evaluation of $`F(x;\lambda )`$ at the points $`\lambda =\lambda _j`$, i.e.,
$$\mu _j=F(x;\lambda _j),$$
provides a set of $`n`$ remaining DN coordinates.
In the application of Section 4, we will use the property that if $`F(x;\lambda )`$ is an exact eigenvector, then $`Y(F(x;\lambda ))`$ is an exact eigenvector as well. Since in the separable case a suitable combination of the Hamiltonians is exact, one can act with $`Y`$ on such a combination and generate a space of exact eigenvectors where the equation $`Y(F(x;\lambda ))=1`$ may be solved algebraically.
### 2.1 DN separable Hamiltonians from GZ systems
Let $``$ a $`(2n+k)`$-dimensional manifold endowed with a pencil $`P_\lambda =P_1\lambda P_0`$ of Poisson tensors. We suppose that it admits $`k`$ polynomial Casimir functions
$$H^{(a)}=\underset{j=0}{\overset{n_a}{}}H_j^{(a)}\lambda ^{n_aj},a=1,\mathrm{},k,$$
with $`n=n_1+\mathrm{}+n_k`$. If the functions $`H_j^{(i)}`$ are functionally independent, then $``$ is called a complete GZ manifold, and the pencil $`P_\lambda `$ is said to be a pure Kronecker pencil of type $`\{2n_1+1,\mathrm{},2n_k+1\}`$. Since the functions $`H_0^{(a)}`$ form a maximal set of independent Casimirs of $`P_0`$, the generic symplectic leaf $`𝒮`$ of $`P_0`$ is the $`2n`$–dimensional submanifold given by $`H_0^{(a)}=C_a`$, for $`a=1,\mathrm{},k`$. The restrictions $`\widehat{H}_{j_a}^{(a)}`$ to $`𝒮`$, for $`j_a=1,\mathrm{},n_a`$, and $`a=1,\mathrm{},k`$, of the $`n`$ remaining Hamiltonians define a completely integrable system in the Liouville sense.
In order to solve by SoV this system, we suppose that there exist $`k`$ vector fields $`Z_a`$ (to be called transversal vector fields) spanning a $`k`$–dimensional integrable distribution $`𝒵`$ and satisfying:
a) The normalized transversality condition: $`\text{Lie}_{Z_a}(H_0^{(b)})=\delta _a^b`$ for all $`a,b=1,\mathrm{},k`$;
b) The deformation condition for the Lie derivatives: $`\text{Lie}_{Z_a}(P_\lambda )=_{b=1}^kZ_bY_a^b`$ for some vector fields $`Y_b^a`$;
c) The “flatness” condition: $`\text{Lie}_{Z_a}(\text{Lie}_{Z_b}(H^{(c)}(\lambda )))=0,a,b,c`$.
Conditions a) and b) imply that the distribution $`𝒵`$ is transversal to the symplectic leaves of $`P_0`$, and that the functions vanishing along $`𝒵`$ are a Poisson subalgebra with respect to the Poisson pencil $`P_\lambda `$. Then, as a consequence of the Marsden–Ratiu theorem , we have that:
###### Proposition 2.5
The Poisson pencil on $``$ can be projected on the generic symplectic leaf $`𝒮`$ of $`P_0`$, so that $`𝒮`$ becomes a PN manifold. The functions $`\widehat{H}_{j_a}^{(a)}`$, for $`j_a=1,\mathrm{},n_a`$, and $`a=1,\mathrm{},k`$, satisfy the condition (2.2), that is, there exists an $`n\times n`$ matrix $`𝖥`$ such that $`N^{}d\widehat{H}=𝖥\widehat{H}`$, where $`\widehat{H}`$ is a column vector collecting the above functions.
Under the “flatness” condition c), one can show that equation (2.4) is satisfied, so that the reduced Hamiltonian system is separable in the DN coordinates. These coordinates may be computed from the geometry of $`P_\lambda `$, without actually performing the reduction process. In this case, in fact:
1. The minimal polynomial of the Nijenhuis tensor $`N`$ induced, according to the previous proposition, on the leaf $`𝒮`$ is the determinant of the matrix $`G(\lambda )=\left[\text{Lie}_{Z_a}(H^{(b)}(\lambda ))\right]_{a,b=1,\mathrm{},k}`$, that is, $`detG(\lambda )=0`$ iff $`\lambda =\lambda _j`$;
2. The vector field $`Y`$ of Theorem 2.4 is given by $`Y=_{a=1}^kY_a^a`$;
3. If $`(1,\rho _2(\lambda ),\mathrm{},\rho _k(\lambda ))`$ satisfies
$$(1,\rho _2(\lambda ),\mathrm{},\rho _k(\lambda ))G(\lambda )=0\text{for }\lambda =\lambda _j\text{,}$$
then $`H^{(1)}(\lambda )+\rho _2(\lambda )H^{(2)}(\lambda )+\mathrm{}+\rho _k(\lambda )H^{(k)}(\lambda )`$ is an exact eigenvector of $`N^{}`$. Hence it can be used to find a normalized exact eigenvector, and therefore the $`\mu _j`$ coordinates.
###### Remark 2.6
The SoV theory for PN manifolds outlined above provides intrinsic and algorithmic recipes to check whether a given Liouville integrable system defined on a PN manifold can be separated in the DN coordinates. On the other hand, the conditions under which one obtains separable Hamiltonians from a GZ manifold are by no means algorithmic. In particular, the existence of the distribution $`𝒵`$ (that is, of the vector fields $`Z_a`$ and $`Y_a^b`$ fulfilling the above three properties) must be checked (and guessed) case by case. In some GZ systems, obtained from stationary reductions of the Boussinesq and KdV hierarchies, are discussed along these lines. In the next sections we will apply the scheme herewith outlined to the Toda lattices.
## 3 The Bihamiltonian approach to Toda lattices
The phase space of the (complex, periodic) Toda lattice (see, e.g., ) with $`n`$ sites (particles) is the manifold $`=(^{})^n\times ^n`$ parametrized by the Flaschka coordinates $`\{a_i,b_i\}_{i=1,\mathrm{},n}`$. We endow it with the Poisson pencil $`P_\lambda `$ defined as follows (see, e.g., and references cited therein). It associates with the one–form $`_k(\alpha _kda_k+\beta _kdb_k)`$ the vector field $`_k(\dot{a}_k_{a_k}+\dot{b}_k_{b_k})`$ according to the rule
$$\begin{array}{c}\dot{a_k}=a_k((b_k\lambda )\beta _k(b_{k+1}\lambda )\beta _{k+1}+a_{k1}\alpha _{k1}\alpha _{k+1}a_{k+1})\hfill \\ \dot{b_k}=(b_k\lambda )(a_{k1}\alpha _{k1}a_k\alpha _k)+a_k\beta _{k+1}a_{k1}\beta _{k1}\hfill \end{array}$$
(3.1)
where the cyclicity condition $`()_{k+n}=()_k`$ is implicitly assumed. We write the matrix expression of $`P_\lambda =P_1\lambda P_0`$ in the $`3`$–particle case, the $`n`$–particle case being easily generalized from this example:
$$P_\lambda =\left[\begin{array}{cccccc}0& a_1a_2& a_1a_3& a_1(b_1\lambda )& a_1(b_2\lambda )& 0\\ \multicolumn{6}{c}{}\\ & 0& a_2a_3& 0& a_2(b_2\lambda )& a_2(b_3\lambda )\\ \multicolumn{6}{c}{}\\ & & 0& a_3(b_1\lambda )& 0& a_3(b_3\lambda )\\ \multicolumn{6}{c}{}\\ & & & 0& a_1& a_3\\ \multicolumn{6}{c}{}\\ & & & & 0& a_2\\ \multicolumn{6}{c}{}\\ & & & & & 0\end{array}\right]$$
(3.2)
According to the GZ scheme, we study the kernel of $`P_\lambda `$. We have to solve the equations
$$\begin{array}{c}(b_k\lambda )\beta _k(b_{k+1}\lambda )\beta _{k+1}+a_{k1}\alpha _{k1}a_{k+1}\alpha _{k+1}=0\hfill \\ (b_k\lambda )(a_{k1}\alpha _{k1}a_k\alpha _k)+a_k\beta _{k+1}a_{k1}\beta _{k1}=0\hfill \end{array}$$
With algebraic manipulations (see ), it can be traded for the system of equations
$$\begin{array}{c}(b_k\lambda )\beta _k+a_{k1}\alpha _{k1}+a_k\alpha _k=L_1\hfill \\ (a_k\alpha _k)^2+a_k\beta _k\beta _{k+1}L_1\alpha _k=L_2\hfill \end{array}$$
(3.3)
where $`L_i`$ are $`_n`$–invariant functions. Setting $`L_1=1,L_2=0`$, and introducing the variables
$$h_k=\frac{\beta _{k+1}}{\alpha _k},$$
we obtain the following Riccati type equation:
$$h_kh_{k+1}=(b_{k+1}\lambda )h_k+a_k.$$
(3.4)
###### Proposition 3.1
The characteristic equation (3.4) admits a solution $`h_k`$ which is a Laurent series in the parameter $`\lambda `$ of the form $`h_k=\lambda +_{j=1}^{\mathrm{}}h_{k,j}\lambda ^j`$. The Laurent coefficients $`h_{k,j}`$ can be computed by recurrence as functions of the variables $`\{a_i,b_i\}`$. The product $`C=h_1\mathrm{}h_n`$ of the components of any solution of (3.4) is a Casimir function of the Poisson pencil $`P_\lambda `$.
Notice that, once the characteristic equation is solved, the one–forms in the kernel of $`P_\lambda `$ can be easily computed (by recurrence) solving the system
$$\{\begin{array}{c}h_k\alpha _k+a_k\beta _k=1\hfill \\ \alpha _kh_k=\beta _{k+1}\hfill \end{array}k=1,\mathrm{},n,$$
which is equivalent to the system (3.3) with $`L_1=1,L_2=0`$.
This method allows us to find Casimirs of $`P_\lambda `$ that are Laurent series in $`\lambda `$. According to the GZ scheme , however, we should better look for polynomial Casimirs of $`P_\lambda `$. They can be found linearizing the Riccati equation (3.4) as follows.
Setting $`h_k=\mu \psi _k/\psi _{k1}`$, we transform equation (3.4) into the linear system
$$\mu ^2\psi _{k+1}\mu (b_k\lambda )\psi _ka_k\psi _{k1}=0,$$
(3.5)
where $`\mu `$ is related to the Casimir $`C`$ via $`C=\mu ^n`$. In matrix form we have $`\psi =0`$, where
$$=\left[\begin{array}{ccccc}\mu (b_1\lambda )& \mu ^2& 0& & a_n\\ a_1& \mu (b_2\lambda )& \mu ^2& \mathrm{}& \\ 0& a_2& \mathrm{}& \mathrm{}& 0\\ & \mathrm{}& \mathrm{}& \mu (b_{n1}\lambda )& \mu ^2\\ \mu ^2& 0& & a_{n1}& \mu (b_n\lambda )\end{array}\right].$$
(3.6)
This is how the classical Lax matrix of the Toda lattice can be introduced into the game in the GZ bihamiltonian point of view. We remark that, since the Riccati equation (3.4) admits solutions, so does the linear system $`\psi =0`$. So, taking into account the cyclicity of $``$, we arrive at
###### Proposition 3.2
The spectral curve of the problem, $`det()=0`$, is a quadratic polynomial in the Casimir $`C`$,
$$det()=C^2+H^{(1)}(\lambda )C+H^{(2)}.$$
(3.7)
Thus, both $`H^{(1)}(\lambda )`$ and $`H^{(2)}`$ are polynomial Casimirs of $`P_\lambda `$. In particular,
$$H^{(2)}=(1)^{n+1}a_1\mathrm{}a_n$$
is a common Casimir of $`P_1`$ and $`P_0`$, and $`H^{(1)}(\lambda )`$ has the form
$$H^{(1)}(\lambda )=(1)^n\lambda ^n+\underset{j=1}{\overset{n}{}}H_j^{(1)}\lambda ^{nj}.$$
It can be easily realized that $`H_j^{(1)}=(1)^j\sigma _j^n(b_1,\mathrm{},b_n)`$+lower order terms in the $`b_j`$, where $`\sigma _j^n`$ is the $`j`$–th elementary symmetric polynomial in $`n`$ letters. So, the Hamiltonian functions $`(H^{(2)},\{H_j^{(1)}\}_{j=1,\mathrm{},n})`$ are functionally independent and the previous proposition provides another proof of the fact that the periodic Toda lattice with $`n`$ particles is a complete GZ manifold of type $`\{1,2n1\}`$.
We end this section with a remark which frames the open Toda lattice within this scheme. It is well known that the open Toda lattice can be obtained form the periodic one by pulling one particle to infinity, that is, in the Flaschka coordinates, by letting one of the $`a`$ coordinates, say $`a_n`$, attain the $`0`$ value. The phase space of the (complex) open Toda lattice with $`n`$ particles is thus the manifold $`\widehat{}=(^{})^{n1}\times ^n`$ parametrized by reduced Flaschka coordinates $`\{a_1\mathrm{},a_{n1},b_1,\mathrm{},b_n\}`$. The Poisson pencil $`\widehat{P_\lambda }`$ of the open case can be obtained from the periodic one by means of the following trick. Let $`\stackrel{~}{}`$ be the manifold obtained from the phase space of the periodic lattice adjoining the “boundary component” defined by $`a_n=0`$. The Poisson pencil defined by (3.1), being polynomial in the Flaschka coordinates, extends naturally to a Poisson pencil $`\stackrel{~}{P}_\lambda `$ on the extended manifold $`\stackrel{~}{}`$. The phase space $`\widehat{}`$ of the open case can be identified with the zero set of the common Casimir $`H^{(2)}`$ of $`\stackrel{~}{P}_1`$ and $`\stackrel{~}{P_0}`$, which, obviously enough, is still given by $`H^{(2)}=(1)^{n+1}a_1\mathrm{}a_n`$. Then $`\stackrel{~}{P}_\lambda `$ can be restricted to $`\widehat{}`$, and its restriction is the Poisson pencil of the open Toda lattice. In practice, its matrix representation in the reduced coordinates is obtained by the matrix representation of the periodic Poisson pencil (3.2) deleting the $`n`$–th row and column, and setting $`a_n=0`$ in the resulting matrix. The Lax matrix of the open Toda Lattice (as well as the Hamiltonian functions) is obtained simply by setting $`a_n=0`$ in the Lax matrix (3.6) of the periodic problem. In particular, the single polynomial Casimir of the Poisson pencil of the open lattice is obtained as $`\widehat{H}^{(1)}=H^{(1)}{}_{|_{a_n=0}}{}^{}.`$ The open Toda lattice is thus a complete GZ manifold of type $`2n1`$.
## 4 Separation of variables
In this Section we will show that the Toda lattice fits the scheme described in Subsection 2.1. We will follow the path of Section 3, considering at first the periodic lattice, and then stating the suitable changes to be done in the open case.
The periodic Toda lattice is a GZ manifold of dimension $`2n`$ and type $`\{1,2n1\}`$. Thus the rank of the transversal distribution $`𝒵`$ must be $`2`$, and the dimension of the reduced PN manifold $`2n2`$. We divide the procedure outlined in Subsection 2.1 in three steps.
Step 1. The transversal vector fields $`Z_1`$ and $`Z_2`$.
###### Proposition 4.1
The vector fields $`Z_1=_{b_n}`$ and $`Z_2=_{a_n}/(a_1\mathrm{}a_{n1})`$ satisfy
$$\text{Lie}_{Z_1}P_\lambda =Z_1Y_{1,1},\text{Lie}_{Z_2}P_\lambda =Z_1Y_{2,1},$$
(4.1)
with $`Y_{1,1}=a_{n1}_{a_{n1}}a_n_{a_n}`$ and $`Y_{2,1}=_{b_1}/(a_1\mathrm{}a_{n1})`$.
This property is proven making use of the standard formulas for the Lie derivative of a bivector.
Step 2. The action of $`Z_i`$ on the Casimirs and the $`\lambda _j`$ coordinates.
To discuss this issue it is useful to recall the expression of the second Casimir $`H^{(2)}=(1)^{n+1}a_1\mathrm{}a_n`$ and to expand $`\text{det}()`$ with respect to the last column to get:
$$\text{det}()=\mu (b_n\lambda )\widehat{L}_{n,n}+\mu ^2\widehat{L}_{n,n1}+(1)^{n+1}a_n\widehat{L}_{n,1},$$
(4.2)
where $`\widehat{L}_{i,j}`$ are the determinants of the suitable minors of $``$. Taking into account the specific form of these minors one can easily see that it holds
###### Proposition 4.2
The second Lie derivatives of $`H^{(1)}`$ and $`H^{(2)}`$ with respect to $`Z_i`$ vanish. Furthermore, $`\text{Lie}_{Z_1}(H^{(2)})=0`$ and $`\text{Lie}_{Z_2}(H^{(2)})=1`$, so that the matrix $`G_a^b=\text{Lie}_{Z_a}(H^{(b)}(\lambda ))`$ introduced in Subsection 2.1 has the form
$$G(\lambda )=\left(\begin{array}{cc}\text{Lie}_{Z_1}(H^{(1)})& 0\\ \text{Lie}_{Z_2}(H^{(1)})& 1\end{array}\right).$$
(4.3)
Thus the $`\lambda _j`$ coordinates are the roots of the monic degree $`(n1)`$ polynomial $`\text{Lie}_{Z_1}(H^{(1)})=_{b_n}det()=\widehat{L}_{n,n}`$.
Step 3. The action of the vector field $`Y`$ and the $`\mu _j`$ coordinates.
We have to consider the vector field $`Y=Y_{1,1}=a_{n1}_{a_{n1}}a_n_{a_n}`$, and to discuss its action on the exact eigenvectors of $`N^{}`$. According to the discussion following Proposition 2.5, to construct such an eigenvector we must find a vector $`(1,\rho (\lambda ))`$ such that $`(1,\rho (\lambda ))G(\lambda )=0`$ for $`\lambda =\lambda _j`$. Since this vector is simply given by $`(1,0)`$, we have that $`H^{(1)}`$ is an exact eigenvector of $`N^{}`$, and this is true, for all $`r`$, for $`Y^r(H^{(1)})`$ as well. In order to build a normalized exact eigenvector, we have to analyze a bit further the terms in the expansion (4.2) of the determinant of $``$. Actually, one has that:
1. $`\widehat{L}_{n,n}`$ is independent of $`a_n`$ and $`a_{n1}`$;
2. $`a_n\widehat{L}_{n,1}=H^{(2)}+CK_1`$ where $`K_1`$ is linear in $`a_n`$ and does not depend on $`a_{n1}`$ and $`\mu `$;
3. $`\mu ^2\widehat{L}_{n,n1}=C^2+CK_2`$ where $`K_2`$ is linear in $`a_{n1}`$ and does not depend on $`a_n`$ and $`\mu `$.
Thanks to the linearity properties of $`K_j`$, we have that $`Y(H^{(1)})=Y(K_1+K_2)`$ satisfies the recursion property $`Y^3(H^{(1)})=Y(H^{(1)})`$. This ensures that the function
$$F=\mathrm{log}(Y(H^{(1)})+Y^2(H^{(1)}))$$
satisfies $`Y(F)`$=1, and, according to Theorem 2.4, is the desired generator of the $`\mu _j`$ coordinates. We notice that, due to the cyclic nature of the periodic Toda system, the pair $`Z_1,Z_2`$ of deformation vector fields is by no means unique; other admissible pairs can be obtained via a cyclic permutation of the indices (with a corresponding change in the vector field $`Y`$). It would be interesting to compare the DN coordinates considered here with the separation variables for Toda systems used in, e.g., .
Finally, we state the corresponding results for the open Toda lattice. We have to look for a single deformation vector field, say $`Z`$; it is still given by $`_{b_n}`$; the vector field $`Y`$ now is given by $`a_{n1}_{a_{n1}}`$; the recursion relation on $`Y(\widehat{H}^{(1)})`$ closes at the first step, $`Y^2(\widehat{H}^{(1)})=Y(\widehat{H}^{(1)})`$, and the generating function can be taken as $`\widehat{F}=\mathrm{log}\left(Y(\widehat{H}^{(1)})\right)`$.
## 5 Example: the three–particle case
We close the paper with some explicit expressions for the three–particle case. The Poisson pencil is written in equation (3.2). The Lax matrix is given by
$$=\left(\begin{array}{ccc}\mu \left(b_1\lambda \right)& \mu ^2& a_3\\ \multicolumn{3}{c}{}\\ a_1& \mu \left(b_2\lambda \right)& \mu ^2\\ \multicolumn{3}{c}{}\\ \mu ^2& a_2& \mu \left(b_3\lambda \right)\end{array}\right).$$
(5.1)
The spectral curve is
$$C^2(\lambda ^3+H_0\lambda ^2+H_1\lambda +H_2)C+K=0,$$
with
$`H_0`$ $`=`$ $`(b_1+b_3+b_2),H_1=a_2+b_1b_3+a_1+b_2b_3+a_3+b_1b_2,`$
$`H_2`$ $`=`$ $`(b_1b_2b_3+a_1b_3+b_1a_2+a_3b_2),K=a_1a_2a_3,C=\mu ^3.`$
There are two nontrivial flows, given by:
$$X_1=\{\begin{array}{c}\dot{a}_1=a_1b_2a_1b_1\hfill \\ \dot{b}_1=a_3a_1\hfill \\ \text{and cyclic permutations}\hfill \end{array}X_2=\{\begin{array}{c}\dot{a}_1=a_1b_1b_3+a_1a_3a_1b_2b_3a_1a_2\hfill \\ \dot{b}_1=a_1b_3a_3b_2\hfill \\ \text{and cyclic permutations}\hfill \end{array}$$
(5.2)
The transversal vector fields are $`Z_1=_{b_3}`$ and $`Z_2=_{a_3}/a_1a_2`$, and we have $`Y=a_2_{a_2}a_3_{a_3}`$. The DN coordinates can be found as follows. The roots of the polynomial
$$\text{Lie}_{Z_1}(H^{(1)}(\lambda ))=_{b_3}(H_0\lambda ^2+H_1\lambda +H_2)=\lambda ^2+(b_1+b_2)\lambda (b_1b_2+a_1)$$
are $`\lambda _1`$ and $`\lambda _2`$. Then $`\mu _1`$ and $`\mu _2`$ are given by the function
$$F=\mathrm{log}\left(Y(H^{(1)})+Y^2(H^{(1)})\right)=\mathrm{log}(2a_2\lambda 2a_2b_1)$$
evaluated at $`\lambda =\lambda _1,\lambda _2`$.
For the open case, the Poisson pencil can be computed from (3.2), according to the procedure outlined at the end of Section 3:
$$P_\lambda ^{open}=\left[\begin{array}{ccccc}0& a_1a_2& \left(b_1\lambda \right)a_1& \left(\lambda b_2\right)a_1& 0\\ \multicolumn{5}{c}{}\\ a_1a_2& 0& 0& \left(b_2\lambda \right)a_2& \left(\lambda b_3\right)a_2\\ \multicolumn{5}{c}{}\\ \left(\lambda b_1\right)a_1& 0& 0& a_1& 0\\ \multicolumn{5}{c}{}\\ \left(b_2\lambda \right)a_1& \left(\lambda b_2\right)a_2& a_1& 0& a_2\\ \multicolumn{5}{c}{}\\ 0& \left(b_3\lambda \right)a_2& 0& a_2& 0\end{array}\right],$$
(5.3)
and the spectral curve is the rational curve
$$\begin{array}{cc}\hfill \mu ^3=C& =\lambda ^3+\left(b_2+b_1+b_3\right)\lambda ^2\left(b_1b_2+a_1+b_2b_3+b_1b_3+a_2\right)\lambda \hfill \\ & +b_1b_2b_3+a_1b_3+b_1a_2.\hfill \end{array}$$
The separation variables can be constructed, mutatis mutandis, as in the periodic case.
### Acknowledgments
G.F. and M.P. would like to thank the organizers of NEEDS99 for the opportunity to present these results there, and for financial support to their participation. This work was partially supported by the M.U.R.S.T. and the G.N.F.M. of the Italian C.N.R. |
warning/0002/hep-ph0002183.html | ar5iv | text | # 1 Couplings and Decay of the Radion to the SM particles
## 1 Couplings and Decay of the Radion to the SM particles
Recently Goldberger and Wise showed that the modulus in the Randall-Sundrum scenario can be stabilized by introducing a scalar field in the bulk. The stabilized modulus turns out to be very light if its mass arises from a small bulk scalar mass. It was shown subsequently that the radion couples to the SM fields on the visible brane via the trace of the energy momentum tensor. In this article we shall study the production and decay of such a light radion at hadron colliders.
The radion couples to the SM particles on the visible brane as given in via the relation, $`_{int}=\frac{1}{\varphi }T_\mu ^\mu \stackrel{~}{\varphi }`$ where $`T_\mu ^\mu `$ is the trace of the symmetrized and conserved energy momentum tensor for SM fields. At the tree level it is given by,
$$T_\mu ^\mu =\underset{f}{}m_f\overline{f}f+2M_z^2Z^\mu Z_\mu +2M_w^2W^\mu W_\mu +2m_h^2h^2.$$
(1)
The fermion and gauge boson terms show the scale breaking effects due to electroweak symmetry breaking. Since the gluons and photons are massless the radion does not couple to these at the tree level. However the running of the gauge coupling in QCD and QED breaks the scale invariance and induces a trace anomaly . The trace anomaly in QCD therefore generates the radion coupling to gluons which is given by
$$_{\stackrel{~}{\varphi }gg}=\frac{1}{\varphi }\frac{\beta (g_s)}{2g_s}\stackrel{~}{\varphi }G^{a\mu \nu }G_{\mu \nu }^a.$$
(2)
where $`\frac{\beta (g_s)}{2g_s}=(11\frac{2}{3}n_f)\frac{g_s^2}{32\pi ^2}`$. $`G_{\mu \nu }^a=[_\mu g_\nu ^a_\nu g_\mu ^a+gf^{abc}g_\mu ^bg_\nu ^c]`$ is the gluon field strength tensor. For $`m_{\stackrel{~}{\varphi }}^2<4m_t^2`$ we have $`n_f=5`$ dynamical quarks and hence $`\frac{\beta (g_s)}{2g_s}3.84\frac{\alpha _s}{4\pi }`$. On the other hand for $`m_{\stackrel{~}{\varphi }}^2>4m_t^2`$ we have $`n_f=6`$ and hence $`\frac{\beta (g_s)}{g_s}3.50\frac{\alpha _s}{4\pi }`$. The anomaly contribution is independent of fermion mass. Even if EW symmetry were exact and all fermions had remained massless the trace anomaly would still lead to the above coupling of $`\stackrel{~}{\varphi }`$ to gluons. In the presence of EWSB, heavy quark loops give rise to a contribution to $`_{\stackrel{~}{\varphi }gg}`$. In the infinite mass limit this contribution can be obtained by replacing $`v`$ by $`\varphi `$ in the effective Lagrangian for $`hgg`$ coupling \[eqn. 10\]. It can be shown that this contribution is smaller than the anomaly contribution written above. The running of QED coupling also introduces a conformal anomaly. This gives rise to the following coupling of the radion to the photons.
$$_{\stackrel{~}{\varphi }\gamma \gamma }=\frac{1}{\varphi }\frac{\beta (e)}{2e}\stackrel{~}{\varphi }F_{\mu \nu }F^{\mu \nu }$$
(3)
Where $`F^{\mu \nu }`$ has the usual meaning. We have used the following values of $`\frac{\beta (e)}{2e}`$ in our calculation.
$`{\displaystyle \frac{\beta (e)}{2e}}`$ $`=`$ $`{\displaystyle \frac{13\alpha }{12\pi }}m_{\stackrel{~}{\varphi }}>2m_t`$ (4)
$`=`$ $`{\displaystyle \frac{31\alpha }{36\pi }}2m_t>m_{\stackrel{~}{\varphi }}>2m_W`$
$`=`$ $`{\displaystyle \frac{10\alpha }{9\pi }}m_{\stackrel{~}{\varphi }}<2m_W`$
We want to mention that in the SM, in the heavy quark limit, the $`h\gamma \gamma `$ coupling is given by,
$$_{h\gamma \gamma }=\frac{1}{v}\frac{\beta (e)}{2e}hF_{\mu \nu }F^{\mu \nu }$$
(5)
where,
$$\frac{\beta (e)}{2e}=\frac{2\alpha }{9\pi }$$
Hence for $`m_{\stackrel{~}{\varphi }}^2,m_h^2<4m_w^2`$ (where $`\stackrel{~}{\varphi }\gamma \gamma `$ is significant and the heavy top limit is valid) we have $`\frac{g_{\stackrel{~}{\varphi }\gamma \gamma }}{g_{h\gamma \gamma }}=5\frac{v}{\varphi }`$. Although $`\stackrel{~}{\varphi }f\overline{f}`$ and $`\stackrel{~}{\varphi }VV`$ ($`V=W,Z`$) couplings are suppressed relative to $`hf\overline{f}`$ and $`hVV`$ couplings for $`\varphi =`$ 1 TeV, $`\stackrel{~}{\varphi }gg`$ and $`\stackrel{~}{\varphi }\gamma \gamma `$ couplings are slightly enhanced relative to $`hgg`$ and $`h\gamma \gamma `$ couplings.
Now we are in a position to present the radion decay branching ratios into different channels. In fig. 1, we present the relevant decay branching ratios. For the illustrative purpose we assumed $`\varphi =`$ 1 TeV. The radion coupling to any SM field has the $`\frac{1}{\varphi }`$ dependence. Thus if we choose a different value of $`\varphi `$, the partial decay widths of different channels change in the same fashion. But the branching ratios remain unchanged. The striking feature is that, for some radion masses the gluon gluon branching ratio is almost equal to 1. The other thing to note is that the photon photon branching ratio is also 5 to 8 times larger than the 2 photon branching ratio of the SM higgs. We see in the next section that this will have some interesting consequence in radion search. In this plot the SM higgs mass of 120 GeV is used for the purpose of illustration. This value of the SM higgs is well above the current experimental lower bound from the LEP . Another interesting feature of this plot is that the radion branching ratio into gg or $`\gamma \gamma `$ increases slowly with $`m_{\stackrel{~}{\varphi }}`$ after the sharp fall around WW threshold. On the contrary the higgs branching ratio in these two modes decreases with $`m_h`$ after the WW threshold.
## 2 Radion Production at Hadron Colliders
The radion coupling to weak gauge bosons are suppressed relative to $`hWW`$ and $`hZZ`$ couplings by a factor of $`\frac{v}{\varphi }`$. It is known that higgs boson production at LHC by the weak boson fusion mechanism is itself suppressed relative to the gluon fusion mechanism over most of the range of $`M_h`$ ($`M_h<`$ 1 TeV). Hence the radion production at LHC by weak boson fusion is also not expected to be the dominant or efficient mechanism. Also the radion couplings to light valence quarks are extremely small. Therefore in this article we shall focus on radion production by gluon fusion. The radion production cross-section at a hadron collider via gluon fusion mechanism is given by,
$$\sigma (pp(\overline{p})\stackrel{~}{\varphi }+X)=\mathrm{\Gamma }(\stackrel{~}{\varphi }gg)\frac{\pi ^2}{8m_{\stackrel{~}{\varphi }}^3}\tau _\tau ^1\frac{dx}{x}g(x)g(\frac{\tau }{x})$$
(6)
Here $`\tau `$ is a dimensionless variable given by $`\frac{m_{\stackrel{~}{\varphi }}^2}{S}`$, where $`S`$ is the proton - proton (anti-proton) center of mass energy. Note that $`L_{\varphi gg}`$ generates momentum dependent $`\stackrel{~}{\varphi }gg`$, $`\stackrel{~}{\varphi }ggg`$ and $`\stackrel{~}{\varphi }ggg`$ couplings. The strength of these couplings are proportional to $`\beta (g_s)`$ which includes the contribution of gluons as well as the dynamical quarks. The Lagrangian also indicates that radions could be produced either singly or in association with gluons at a hadron collider by gluon fusion. In fact the gluon fusion mechanism turns out to be the dominant production process for radions at hadron colliders over most of the interesting range. $`qg\stackrel{~}{\varphi }`$ and $`\overline{q}g\stackrel{~}{\varphi }`$ does make a small contribution to radion production in $`O(\alpha _s)^3`$. The radion coupling to gluons is very similar in structure to the effective Lagrangian that gives the higgs coupling to gluons in the heavy quark limit. As in the case of the radion the gluon fusion also turns out to be the primary production mechanism for higgs boson at hadron collider. The dominant contribution to $`ggh`$ arises from closed loops of heavy quarks that occur in the theory. In this work we shall assume that the number of heavy quarks ($`N_h`$) is equal to one namely the top. It has been shown that the heavy quark limit $`\frac{M_h}{M_q}0`$ is an excellent approximation to the exact two loop corrected rate for $`ggh`$. As $`m_h2m_t`$ the exact result rises above the heavy quark result and exhibits a small bump corresponding to the $`t\overline{t}`$ threshold. The width of the bump i.e the departure region increases with increasing $`m_t`$. The disagreement between the two results in the $`m_h2m_t`$ region however is always less than a factor of two at LHC. The heavy quark limit for $`ggh`$ can be obtained from the gauge invariant effective Lagrangian
$$=\frac{1}{4}\left[1\frac{2\beta _h}{g_s(1+\delta )}\frac{h}{v}\right]G_{\mu \nu }^aG^{a\mu \nu }\frac{m_t}{v}h\overline{t}t.$$
(7)
$`\delta =1+2\frac{\alpha _s}{\pi }`$ is the anomalous dimension of the mass operator arising from QCD interactions. $`\beta _h`$ is the heavy quark contribution to the QCD beta function. Since the $`hgg`$ coupling in the $`M_q\mathrm{}`$ limit arises from heavy quark loops it is only the heavy quarks that contribute to the $`\beta _h`$ in eqn(3). To order $`(\alpha _s^3)`$ the heavy quark contribution to $`\beta (g_s)`$ is given by $`\beta _h=N_h\frac{\alpha _s}{12\pi }[1+\frac{19\alpha _s}{4}]`$. On the other hand the $`\beta (g_s)`$ that appears in the gluon coupling to the dilaton like radion mode arises from the trace anomaly. The trace anomaly has its origin in the heavy regulator fields as their masses are taken to infinity. So it includes the gluonic contribution as well as that of dynamical quarks. This difference in the two beta function contributions makes the $`\stackrel{~}{\varphi }gg`$ coupling greater than the $`hgg`$ coupling even for $`\frac{v}{\varphi }=\frac{1}{4}`$. However with increasing $`\varphi `$ the $`hgg`$ coupling ultimately wins over the $`\stackrel{~}{\varphi }gg`$ coupling. This feature is clear from fig. 2 where it is shown that with increasing $`\varphi `$, $`\sigma (pp\stackrel{~}{\varphi })`$ ultimately becomes smaller than $`\sigma (pph)`$.
Let us now make some rough numerical estimate about the ratio $`\frac{\sigma (pp\varphi )}{\sigma (pph)}`$ in the lowest order (O($`\alpha _s`$)). Our estimates will depend only on the relative strength of $`\varphi gg`$ and $`hgg`$ couplings. For $`\sqrt{\widehat{s}}<2m_t`$, $`\frac{\beta (g_s)}{2g_s}=3.84\frac{\alpha _s}{4\pi }`$ whereas $`\beta _h=\frac{1}{3}\frac{\alpha _s}{4\pi }`$ to lowest order. Also in this region the heavy quark limit provides a good approximation to the exact result for $`\sigma (pph)`$. We find that $`\frac{g_{gg\varphi }}{g_{ggh}}2.88`$ for $`\varphi =`$ 1 TeV. Above the $`2m_t`$ threshold we have $`\frac{\beta (g_s)}{2g_s}3.50\frac{\alpha _s}{4\pi }`$. Although in this region the heavy quark limit does not work that well for the higgs cross section we can still get an order of magnitude estimate (lower by at most a factor of two) using it. The ratio of couplings now ($`\sqrt{\widehat{s}}>2m_t`$) becomes $`\frac{g_{gg\varphi }}{g_{ggh}}2.63`$.
Using the fact that the effective Lagrangian for higgs and radion production by gluon fusion are similar except for couplings we find that if $`\varphi `$=1 TeV then the cross section for $`pp\stackrel{~}{\varphi }`$ will exceed that of $`pph`$ by a factor of 8.3 for $`\sqrt{\widehat{s}}<2m_t`$ and by a factor of 6.9 for $`\sqrt{\widehat{s}}>2m_t`$. However for $`\varphi `$=5 TeV the radion production cross section will be suppressed relative to the higgs cross section roughly by a factor of three. These features have been exhibited in fig. 2 where we have plotted the lowest order higgs production cross section (both exact and the heavy quark limit) and the radion cross section (for three different values of $`\varphi `$) against the mass of the particle ($`\stackrel{~}{\varphi }`$ or h) at LHC.
The estimates given above are based on lowest order calculations. It is known that higher order QCD corrections increases the lowest order rate by a factor (K factor) that lies between 2 and 3 at LHC . The QCD radiative corrections done in the heavy quark limit forms an excellent approximation to the exact calculations. To calculate the K factor one therefore always uses the heavy quark limit. But in the heavy quark limit the effective Lagrangian for higgs production is similar in structure to the Lagrangian for radion production. Hence the K factor for higgs production in the heavy quark limit will be the same for the radion also. So higher order QCD corrections will not affect the relative rate between the radion and the higgs to a very high degree of accuracy.
## 3 Detection and Possible SM Backgrounds
Let us now concentrate on the detection of radion at a hadron collider like Tevatron. The dominant decay mode of a 50-150 GeV radion as can be seen from fig. 1, is to $`b\overline{b}`$ or to $`gg`$. But the striking feature is that the $`\gamma \gamma `$ branching ratio of the radion is larger than the higgs case by a factor 5-8, over a considerable mass range. The higgs production rate by gluon fusion and and its decay into the $`\gamma \gamma `$ mode are both suppressed relative to that of the radion. This is the reason why two photon final state is not a good bet for the higgs at the Tevatron.
At the Tevatron radion production cross-section varies from 140 $`pb`$ to 1 $`pb`$ as we vary radion mass from 20 (current lower bound on radion mass comes from LEP -II ) to 160 GeV. We have included a NLO QCD correction factor of 2 in our calculation. This cross-section with the presently collected luminosity will give rise to some 10 <sup>4</sup> radions. If the radion decays into the $`\gamma \gamma `$ mode the final state will consist of 2 hard photons. The main background of this 2-photon final state comes from the pair annihilation of the valence quarks and anti-quarks. The other dominant source of $`\gamma \gamma `$ background is gluon gluon annihilation to two photons. Though this is suppressed to the former by a factor of $`\alpha _s^2`$, dominance of gluon flux over the quarks flux, can make this comparable with the former. We do not calculate this second contribution explicitly. We multiply the $`q\overline{q}\gamma \gamma `$ contribution by a factor of 2 to take this into account. At the Tevatron this is a conservative approximation. We have used a parton level monte-carlo event generator to estimate the numbers for both the signal and the background and the CTEQ-4M parametrisation for the parton densities in our entire analysis.
The following cuts have been applied to differentiate between signal and background.
$`p_T^\gamma >10`$ GeV
We demand the photons are in the central part of the detector.
$`|\eta _\gamma |<3.`$.
We also require that the angular separation between the photons be substantial, i.e. $`\mathrm{\Delta }R_{jj}\sqrt{(\mathrm{\Delta }\eta _{\gamma \gamma })^2+(\mathrm{\Delta }\varphi _{\gamma \gamma })^2}>0.5`$
Even after applying these cuts, the SM background cannot be removed completely. So the strategy is to compare the invariant mass (of the photon pair) distribution of the signal and background. For the signal (of a definite $`m_{\stackrel{~}{\varphi }}`$), invariant mass distribution shows a sharp peak over the continuum background. The sharpness of the peak depends mainly on the detector resolution and the partial width of $`\stackrel{~}{\varphi }\gamma \gamma `$. In fig. 3 we plot the invariant mass ($`M_{\gamma \gamma }`$) distribution for the signal (dots) and backgrounds (broken histogram) assuming an uniform bin size of 5 GeV. Here one can see that, for low $`M_{\gamma \gamma }`$, number of background events is higher than the signal. But as $`M_{\gamma \gamma }`$ increases, the number of background events in the mass bins falls off more rapidly than the signal- which more or less remains the same over the entire mass range we are interested here. Once the radion mass is near to $`2m_w`$, signal events falls off sharply, due to the sharp fall of $`\gamma \gamma `$ branching ratio.
We find that the $`\gamma \gamma `$ mode is good enough to exclude radion mass nearly upto 120 GeV even with the presently collected luminosity of 110 $`pb^1`$. Here we have taken the radion $`vev`$ to be 1 TeV. On the other hand at the Tevatron with the presently collected luminosity one cannot say anything about the SM higgs.
Now let us go over to the case of Upgraded Tevatron with center of mass energy of 2 TeV and luminosity of 1 $`fb^1`$. As center of mass energy is almost equal to the present, the signal and background are just 10 times larger than the previous case. This is evident from the fig. 4. In this figure we also plot the significance ($`\frac{No.ofSignal}{\sqrt{No.ofBackground}})`$ of our signal upto $`M_{\gamma \gamma }`$ equal to 100 GeV. One can check, over this entire mass range, significance is greater than 5. Generally a significance greater than 5 points to the discovery. One can also easily check that once $`M_{\gamma \gamma }`$ is greater than 100 GeV, no. of background events in the corresponding bins become less than 1 events. So in this region we can demand 5 signal events as a benchmark for discovery. Thus at the upgraded Tevatron one can discover radion mass upto 160 GeV and exclude upto 165 GeV. Once $`M_{\gamma \gamma }`$ is greater than 165 GeV, our signal falls off very sharply and we cannot say anything more about it.
If we change $`\varphi `$ the branching ratio of the radion to different channels remains same. So if we take $`\varphi =`$ 2 TeV, the cross-section and number of $`\gamma \gamma `$ events become $`\frac{1}{4}`$ of the present case ($`\varphi =`$ 1 TeV). And at the Tevatron with the presently collected luminosity, we cannot say anything about it. At the Tevatron Upgrade, we also cannot talk about the discovery, for $`m_{\stackrel{~}{\varphi }}<`$45 GeV. At higher masses, one can discover upto 150 GeV $`m_\varphi `$ if $`\varphi =`$ 2 TeV instead of 1 TeV.
Finally we want to examine the search prospects of this scalar particle at the LHC. Though the two photon branching ratio drops very sharply near $`m_{\stackrel{~}{\varphi }}`$ 140 GeV, in contrast to the SM case it remains constant with $`m_{\stackrel{~}{\varphi }}`$ after 140 GeV. Also the radion production rate is almost 8 times higher than that of higgs production rate. So unlike the SM higgs boson, for $`m_{\stackrel{~}{\varphi }}>`$ 140 GeV, $`\gamma \gamma `$ mode is a viable avenue to discover or exclude the radion for masses well upto 1 TeV.
Conservatively, we take the luminosity to be equal to 30 $`fb^1`$. The source of SM backgrounds remain the same here. But unlike the Tevatron, the 2-photon background coming from pair annihilation of quark and anti-quark becomes less severe. This is because the anti-quark coming from one proton has to be excited from the sea. But gluon density in the proton is much larger than the quark density at LHC. So to take into account the gluon gluon contribution to 2-photon background we multiply the quark-anti-quark contribution by a factor of 6 <sup>1</sup><sup>1</sup>1 Generally at the LHC energies the gluon-gluon contribution to two photon background is 5 times larger than the quark anti-quark contribution. For $`m_{\gamma \gamma }=`$ 110 (130) GeV this factor has been estimated to be 3.55 (3.56) for the LHC .. The cuts we used here are rather similar to the Tevatron case. Only we change the cuts on the transverse momenta of the photons to be greater than 20 GeV. We present in fig. 5 the number of signal events and the corresponding number for the background against the $`\gamma \gamma `$ invariant mass. To plot the invariant mass distribution for the background we choose a bin size of 5 GeV. Form the figure one can easily see that the background is an order of magnitude smaller than the signal over the entire mass range. And once $`M_{\gamma \gamma }`$ is greater than 250 GeV, there are less than one event in the respective mass bins. That’s why we do not show it in the figure. So LHC can easily discover radions with masses upto 1 TeV for $`\varphi `$=1 Tev. If we increase $`\varphi `$ to 4 TeV one can easily check from fig. 5 that our discovery limit comes down to 650 GeV.
## 4 Conclusions
In this paper we examined the radion production and its subsequent decay into SM particles at hadron colliders. The radion production cross-section is larger than the higgs production cross-section by a factor of 6-8. The partial width of the radion to two gluons and two photons is also enhanced due to the enhanced $`\stackrel{~}{\varphi }gg`$ or $`\stackrel{~}{\varphi }\gamma \gamma `$ coupling. The other partial widths are suppressed with respect to higgs widths by a factor that depends on $`\varphi `$. We also discussed the viability of the 2-photon signal for the radion at the Tevatron. One can exclude radion mass upto 120 GeV even with the presently collected luminosity. Upgraded Tevatron with higher center of mass energy and higher luminosity certainly can discover radion upto 160 GeV mass, if not it can exclude it upto 165 GeV. Similarly at the LHC one can definitely discover radions with masses up to 1 TeV. These estimates are done assuming $`\varphi `$ =1 TeV. With increasing $`\varphi `$ both the exclusion limit and the discovery limit using the $`\gamma \gamma `$ mode will come down.
Note added: While this work was being completed there appeared one paper which discusses some of the issues presented here. |
warning/0002/math0002149.html | ar5iv | text | # A combinatorial approach to quantification of Lie algebras
## Abstract.
We propose a notion of a quantum universal enveloping algebra for an arbitrary Lie algebra defined by generators and relations which is based on the quantum Lie operation concept. This enveloping algebra has a PBW basis that admits the Kashiwara crystalization. We describe all skew primitive elements of the quantum universal enveloping algebra for the classical nilpotent algebras of the infinite series defined by the Serre relations and prove that the set of PBW-generators for each of these enveloping algebras coincides with the Lalonde–Ram basis of the ground Lie algebra with a skew commutator in place of the Lie operation. The similar statement is valid for Hall–Shirshov basis of any Lie algebra defined by one relation, but it is not so in general case.
Research at MSRI is supported in part by NSF grant DMS-9701755. Supported in part by CONACyT México, grant 32130-E
1.Introduction
Quantum universal enveloping algebras appeared in the famous papers by Drinfeld and Jimbo . Since then a great deal of articles and number of monographs were devoted to their investigation. All of these researches are mainly concerned with a particular quantification of Lie algebras of the classical series. This is accounted for first by the fact that these Lie algebras have applications and visual interpretations in physical speculations, and then by the fact that a general, and commonly accepted as standard, notion of a quantum universal enveloping algebra is not elaborated yet (see a detailed discussion in ).
In the present paper we propose a combinatorial solution of this problem by means of the quantum (Lie) operation concept . In line with the main idea of our approach, the skew primitive elements must play the same role in quantum enveloping algebras as the primitive elements do in the classical case. By the Friedrichs criteria , the primitive elements form the ground Lie algebra in the classical case. For this reason we consider the space spanned by the skew primitive elements and equipped with the quantum operations as a quantum analogue of a Lie algebra.
In the second section we adduce the main notions and consider some examples. These examples, in particular, show that the Drinfeld–Jimbo enveloping algebra as well as its modifications are quantum enveloping algebras in our sense.
In the third section with the help of the Heyneman–Radford theorem we introduce a notion of a combinatorial rank of a Hopf algebra generated by skew primitive semi-invariants. Then we define the quantum enveloping algebra of an arbitrary rank that slightly generalises the definitions given in the preceding section.
The basis construction problem for the quantum enveloping algebras is considered in the fourth section. We indicate two main methods for the construction of PBW-generators. One of them modifies the Hall–Shirshov basis construction process by means of replacing the Lie operation with a skew commutator. The set of the PBW-generators defined in this way, the values of hard super-letters, plays the same role as the basis of the ground Lie algebra does in the PBW theorem. At first glance it would seem reasonable to consider the $`𝐤[G]`$-module generated by the values of hard super-letters as a quantum Lie algebra. However, this extremely important module falls far short of being uniquely defined. It essentially depends on the ordering of the main generators, their degrees, and it is almost never antipode stable. Also we have to note the following important fact. Our definition of the hard super-letter is not constructive and, of course, it cannot be constructive in general. The basis construction problem includs the word problem for Lie algebras defined by generators and relations, while the latter one has no general algorithmic solution (see ).
The second method is connected with the Kashiwara crystallisation idea (see also a development in ). M. Kashiwara has considered the main parameter $`q`$ of the Drinfeld–Jimbo enveloping algebra as a temperature of some physical medium. When the temperature tend to zero, the medium crystallises. The PBW-generators must crystallise as well. In our case under this process no one limit quantum enveloping algebra appears since the existence conditions normally include equalities of the form $`p_{ij}=1`$ (see ). Nevertheless if we equate all quantification parameters to zero, the hard super-letters would form a new set of PBW-generators for the given quantum universal enveloping algebra. To put this another way, the PBW-basis defined by the super-letters admits the Kashiwara crystallisation.
In the fifth section we bring a way to construct a Groebner–Shirshov relations system for a quantum enveloping algebra. This system is related to the main skew primitive generators, and, according to the Diamond Lemma (see ), it determines the crystal basis. The usefulness of the Groebner–Shirshov systems depends upon the fact that such a system not only defines a basis of an associative algebra, but it also provides a simple diminishing algorithm for expansion of elements on this basis (see, for example ).
In the sixth section we adapt a well known method of triangular splitting to the quantification with constants. The original method appeared in studies of simple finite dimensional Lie algebras. Then it has been extended into the field of quantum algebra in a lot of publications (see, for example ). By means of this method the investigation of the Drinfeld–Jimbo enveloping algebra amounts to a consideration of its positive and negative homogeneous components, quantum Borel sub-algebras.
In the seventh section we consider more thoroughly the quantum universal enveloping algebras of nilpotent algebras of the series $`A_n,`$ $`B_n,`$ $`C_n,`$ $`D_n`$ defined by the Serre relations. We adduce first lists of all hard super-letters in the explicit form, then Groebner–Shirshov relations systems, and next spaces $`L(U_P(𝔤))`$ spanned by the skew primitive elements (i.e. the Lie algebra quantifications $`𝔤_P`$ proper). In all cases the lists of hard super-letters (but the hard super-letters themselves) turn out to be independent of the quantification parameters. This means that the PBW-generators result from the Hall–Shirshov basis of the ground Lie algebra by replacing the Lie operation with the skew commutator. The same is valid for the Groebner–Shirshov relations systems. Note that the Hall–Shirshov bases, under the name standard Lyndon bases, for the classical Lie series were constructed by P. Lalonde and A. Ram , while the Groebner–Shirshov systems of Lie relations were found by L.A. Bokut’ and A.A. Klein .
Furthermore, in all cases $`𝔤_P`$ as a quantum Lie algebra (in our sense) proves to be very simple in structure. Either it is a coloured Lie super-algebra (provided that the main parameter $`p_{11}`$ equals 1), or values of all non-unary quantum operations equal zero on $`𝔤_P.`$ In particular, if char(k) $`=0`$ and $`p_{11}^t1`$ then the partial quantum operations may be defined on $`𝔤_P,`$ but all of them have zero values. Thus, in this case we have a reason to consider $`U_P(𝔤)`$ as an algebra of ‘commutative’ quantum polynomials, since the universal enveloping algebra of a Lie algebra with zero bracket is the algebra of ordinary commutative polynomials. From this standpoint the Drinfeld–Jimbo enveloping algebra is a ‘quantum’ Weyl algebra of (skew) differential operators. Immediately afterwards a number of interesting questions appears. What is the structure of other algebras of ‘commutative’ quantum polynomials? Under what conditions are the quantum universal enveloping algebras of homogeneous components of other Kac–Moody algebras defined by the Gabber–Kac relations the algebras of ‘commutative’ quantum polynomials? When do the PBW-generators result from a basis of the ground Lie algebra by means of replacing the Lie operation with the skew commutator? These and other questions we briefly discuss in the last section.
It is as well to bear in mind that the combinatorial approach is not free from flaws: the quantum universal enveloping algebra essentially depends on a combinatorial representation of the ground Lie algebra, i.e. a close connection with the abstract category of Lie algebras is lost.
2. Quantum enveloping algebras
Recall that a variable $`x`$ is called a quantum variable if an element $`g_x`$ of a fixed Abelian group $`G`$ and a character $`\chi ^xG^{}`$ are associated with it. A noncommutative polynomial in quantum variables is called a quantum operation if all of its values in all Hopf algebras are skew primitive provided that every variable $`x`$ has a value $`x=a`$ such that
(1)
$$\mathrm{\Delta }(a)=a1+g_xa,g^1ag=\chi ^x(g)a,gG.$$
Let $`x_1,\mathrm{},x_n`$ be a set of quantum variables. For each word $`u`$ in $`x_1,\mathrm{},x_n`$ we denote by $`g_u`$ an element of $`G`$ that appears from $`u`$ by replacing of all $`x_i`$ with $`g_{x_i}.`$ In the same way we denote by $`\chi ^u`$ a character that appears from $`u`$ by replacing of all $`x_i`$ with $`\chi ^{x_i}.`$ Thus on the free algebra k$`x_1,\mathrm{},x_n`$ a grading by the group $`G\times G^{}`$ is defined. For each pair of homogeneous elements $`u,`$ $`v`$ we fix the denotations $`p_{uv}=\chi ^u(g_v)=p(u,v).`$
The quantum operation can be defined equivalently as a $`G\times 1`$-homogeneous polynomial that has only primitive values in all braided bigraded Hopf algebras provided that all quantum variables have primitive homogeneous values $`g_a=g_x,`$ $`\chi ^a=\chi ^x`$ (see \[21, Sect. 1–4\]).
Recall that a constitution of a word $`u`$ is a sequence of non-negative integers $`(m_1,m_2,\mathrm{},m_n)`$ such that $`u`$ is of degree $`m_1`$ in $`x_1,`$ deg$`{}_{1}{}^{}(u)=m_1;`$ of degree $`m_2`$ in $`x_2,`$ deg$`{}_{2}{}^{}(u)=m_2;`$ and so on. Since the group $`G`$ is abelian, all constitution homogeneous polynomials are homogeneous with respect to the grading. Let us define a bilinear skew commutator on the set of graded homogeneous noncommutative polynomials by the formula
(2)
$$[u,v]=uvp_{uv}vu.$$
These brackets satisfy the following Jacobi and skew differential identities:
(3)
$$[[u,v],w]=[u,[v,w]]+p_{wv}^1[[u,w],v]+(p_{vw}p_{wv}^1)[u,w]v;$$
(4)
$$[[u,v],w]=[u,[v,w]]+p_{vw}[[u,w],v]+p_{uv}(p_{vw}p_{wv}1)v[u,w];$$
(5)
$$[u,vw]=[u,v]w+p_{uv}v[u,w];[uv,w]=p_{vw}[u,w]v+u[v,w],$$
where by the dot we denote the usual multiplication. It is easy to see that the following conditional restricted identities are valid as well
(6)
$$[u,v^n]=[\mathrm{}[[u,v],v]\mathrm{},v];[v^n,u]=[v,[\mathrm{}[v,u]\mathrm{}]],$$
provided that $`p_{vv}`$ is a primitive $`t`$-th root of unit, and $`n=t`$ or $`n=tl^k`$ in the case of characteristic $`l>0.`$
Suppose that a Lie algebra $`𝔤`$ is defined by the generators $`x_1,\mathrm{},x_n`$ and the relations $`f_i=0.`$ Let us convert the generators into quantum variables. For this associate to them elements of $`G\times G^{}`$ in arbitrary way. Let $`P=p_{ij},`$ $`p_{ij}=\chi ^{x_i}(g_{x_j})`$ be the quantification matrix.
###### Definition 2.1.
A braided quantum enveloping algebra is a braided bigraded Hopf algebra $`U_P^b(𝔤)`$ defined by the variables $`x_1,\mathrm{},x_n`$ and the relations $`f_i=0,`$ where the Lie operation is replaced with (2), provided that in this way $`f_i`$ are converted into the quantum operations $`f_i^{}.`$ The coproduct and the commutation relations in the tensor product are defined by
(7)
$$\mathrm{\Delta }(x_i)=x_i\underset{¯}{}1+1\underset{¯}{}x_i,$$
(8)
$$(x_i\underset{¯}{}x_j)(x_k\underset{¯}{}x_m)=(\chi ^{x_k}(g_{x_j}))^1x_ix_k\underset{¯}{}x_jx_m.$$
###### Definition 2.2.
A simple quantification or a quantum universal enveloping algebra of $`𝔤`$ is an algebra $`U_P(𝔤)`$ that is isomorphic to the skew group algebra
(9)
$$U_P(𝔤)=U_P^b(𝔤)G,$$
where the group action and the coproduct are defined by
(10)
$$g^1x_ig=\chi ^{x_i}(g)x_i,\mathrm{\Delta }(x_i)=x_i1+g_{x_i}x_i,\mathrm{\Delta }(g)=gg.$$
###### Definition 2.3.
A quantification with constants is a simple quantification where additionally some generators $`x_i`$ associated to the trivial character are replaced with the constants $`\alpha _i(1g_{x_i}).`$
The formulae (10) and (7) correctly define the coproduct since by definition of the quantum operation $`\mathrm{\Delta }(f_i^{})=f_i^{}1+gf_i^{}`$ in the case of ordinary Hopf algebras and $`\mathrm{\Delta }(f_i^{})=f_i^{}\underset{¯}{}1+1\underset{¯}{}f_i^{}`$ in the braided case.
We have to note that the defined quantifications essentially depend on the combinatorial representation of the Lie algebra. For example, an additional relation $`[x_1,x_1]=0`$ does not change the Lie algebra. At the same time if $`\chi ^{x_1}(g_1)=1`$ then this relation admits the quantification and yields a nontrivial relation for the quantum enveloping algebra, $`2x_1^2=0.`$
Example 1. Suppose that the Lie algebra is defined by a system of constitution homogeneous relations. If the characters $`\chi ^i`$ are such that $`p_{ij}p_{ji}=1`$ for all $`i,j`$ then the skew commutator itself is a quantum operation. Therefore on replacing the Lie operation all relations become quantum operations as well. This means that the braided enveloping algebra is the universal enveloping algebra $`U(𝔤^{col})`$ of the coloured Lie super-algebra which is defined by the same relations as the given Lie algebra is. The simple quantification appears as the Radford biproduct $`U(𝔤^{col})𝐤[G]`$ or, equivalently, as the universal $`G`$-enveloping algebra of the coloured Lie super-algebra $`𝔤^{col}`$ (see or \[21, Example 1.9\]).
Example 2. Suppose that the Lie algebra $`𝔤`$ is defined by the generators $`x_1,\mathrm{},x_n`$ and the system of nil relations
(11)
$$x_j(adx_i)^{n_{ij}}=0,1ijn.$$
Usually instead of the matrix of degrees (without the main diagonal) $`n_{ij}`$ the matrix $`A=a_{ij},`$ $`a_{ij}=1n_{ij}`$ is considered. The Coxeter graph $`𝚪(A)`$ is associated to every such a matrix. This graph has the vertices $`1,\mathrm{},n,`$ where the vertex $`i`$ is connected by $`a_{ij}a_{ji}`$ edges with the vertex $`j.`$
If $`a_{ij}=0`$ then the relation $`x_j`$ad$`x_i=0`$ is in the list (11), and the relation $`x_i`$(ad$`x_j)^{n_{ji}}=0`$ is a consequence of it. The skew commutator $`[x_j,x_i]`$ is a quantum operation if and only if $`p_{ij}p_{ji}=1.`$ Under this condition we have $`[x_i,x_j]=p_{ij}[x_j,x_i].`$ Therefore both in the given Lie algebra and in its quantification one may replace the relation $`x_i`$(ad$`x_j)^{n_{ji}}=0`$ with $`x_i`$ad$`x_j=0.`$ In other words, without loss of generality, we may suppose that $`a_{ij}=0a_{ji}=0.`$ By the Gabber-Kac theorem we get that the algebra $`𝔤`$ is the positive homogeneous component $`𝔤_1^+`$ of a Kac-Moody algebra $`𝔤_1.`$
Theorem 6.1 describes the conditions for a homogeneous polynomial in two variables which is linear in one of them to be a quantum operation. From this theorem we have the following corollary.
###### Corollary 2.4.
If $`n_{ij}`$ is a simple number or unit and in the former case $`p_{ii}`$ is not a primitive $`n_{ij}`$-th root of unit, then the relation $`(\text{11})`$ admits a quantification if and only if $`p_{ij}p_{ji}=p_{ii}^{a_{ij}}.`$
Theorem 6.1 provides no essential restrictions on the non-diagonal parameters $`p_{ij}:`$ if the matrix $`P`$ correctly defines a quantification of (11) then for every set $`Z=\{z_{ij}|z_{ij}z_{ji}=z_{ii}=1\}`$ the following matrix does as well:
(12)
$$P_Z=\{p_{ij}z_{ij}|p_{ij}P,z_{ij}Z\}.$$
Example 3. Let $`G`$ be freely generated by $`g_1,\mathrm{}g_n`$ and $`A`$ be a generalised Cartan matrix symmetrised by $`d_1,\mathrm{},d_n,`$ while the characters are defined by $`p_{ij}=q^{d_ia_{ij}}.`$ In this case the simple quantification is the positive component of the Drinfeld–Jimbo enveloping algebra together with the group-like elements, $`U_P(𝔤)=U_q^+(𝔤)G.`$ By means of an arbitrary deformation (12) one may define a ‘colouring’ of $`U_q^+(𝔤)G.`$
The braided enveloping algebra equals $`U_q^+(𝔤)`$ where the coproduct and braiding are defined by (7) and (8) with the coefficient $`q^{d_ka_{kj}}.`$ The formula (12) correctly defines its ‘colouring’ as well.
Example 4. If in the example above we complete the set of quantum variables by the new ones $`x_1^{},\mathrm{},x_n^{};`$ $`z_1,\mathrm{},z_n`$ such that
(13)
$$\chi ^x^{}=(\chi ^x)^1,g_x^{}=g_x,\chi ^z=\text{id},g_{z_i}=g_i^2,$$
then by \[21, Theorem 6.1\] the Gabber–Kac relations (2), (3), and $`[e_i,f_j]=\delta _{ij}h_i`$ (see \[16, Theorem 2\]) under the identification $`e_i=x_i,`$ $`f_i=x_i^{},`$ $`h_i=z_i`$ admit the quantification with constants $`z_i=\epsilon _i(1g_i^2).`$ (Unformally we may consider the obtained quantification as one of the Kac–Moody algebra identifying $`g_i`$ with $`q^{h_i},`$ where the rest of the Kac–Moody algebra relations, $`[h_i,e_j]=a_{ij}e_i,`$ $`[h_i,f_j]=a_{ij}f_j,`$ is quantified to the $`G`$-action: $`g_j^1x_i^\pm g_j=q^{d_{ij}a_{ij}}x_i^\pm .`$) This quantification coincides with the Drinfeld–Jimbo one under a suitable choice of $`x_i,`$ $`x_i^{},`$ and $`\epsilon _i`$ depending up the particular definition of $`U_q(𝔤):`$
$$\begin{array}{cc}[\text{28}]& x_i=E_i,g_i=K_i,x_i^{}=F_iK_i,p_{ij}=v^{d_ia_{ij}},\epsilon _i=(v^{d_i}v^{d_i})^1;\\ [\text{29}]& x_i=E_i,g_i=\stackrel{~}{K}_i,x_i^{}=F_i\stackrel{~}{K}_i,p_{i\mu }=v^{\mu ,i^{}},\epsilon _i=(v_i^1v_i)^1;\\ [\text{19}]\mathrm{\Delta }_+& x_i=e_i,g_i=t_i,x_i^{}=t_if_i,p_{ij}=q_j^{h_j,\alpha _i},\epsilon _i=(q_iq_i^3)^1;\\ [\text{19}]\mathrm{\Delta }_{}& x_i=f_i,g_i=t_i,x_i^{}=e_it_i,p_{ij}=q_j^{h_j,\alpha _i},\epsilon _i=(q_i^1q_i)^1;\\ [\text{34}]& x_i=E_iK_i,g_i=K_i^2,x_i^{}=F_iK_i,p_{ij}=q^{2d_ia_{ij}},\epsilon _i=(1q^{4d_i})^1.\end{array}$$
By (13) the brackets $`[x_i,x_j^{}]`$ are quantum operations only if $`p_{ij}=p_{ji}.`$ So in this case the ‘colourings’ (12) may be only black-white, $`z_{ij}=\pm 1.`$
In the perfect analogy the Kang quantification of the generalised Kac-Moody algebras is a quantification in our sense as well.
3.Combinatorial rank
By the above definitions the quantum enveloping algebras (with or without constants) are character Hopf algebras (see \[21, Definition 1.2\]). In this section by means of a combinatorial rank notion we identify the quantum enveloping algebras in the class of character Hopf algebras.
Let $`H`$ be a character Hopf algebra generated by $`a_1,\mathrm{},a_n:`$
(14)
$$\mathrm{\Delta }(a_i)=a_i1+g_{a_i}a_i,g^1a_ig=\chi ^{a_i}(g)a_i,gG.$$
Let us associate a quantum variable $`x_i`$ with the parameters $`(\chi ^{a_i},`$ $`g_{a_i})`$ to $`a_i.`$ Denote by $`GX`$ the free enveloping algebra defined by the quantum variables $`x_1\mathrm{},x_n.`$ (see \[21, Sec. 3\] under denotation $`HX`$).
The map $`x_ia_i`$ has an extension to a homomorphism of Hopf algebras $`\phi :GXH.`$ Denote by $`I`$ the kernel of this homomorphism. If $`I0`$ then by the Heyneman–Radford theorem (see, for example \[34, pages 65–71\]), the Hopf ideal $`I`$ has a non-zero skew primitive element. Let $`I_1`$ be an ideal generated by all skew primitive elements of $`I.`$ Clearly $`I_1`$ is a Hopf ideal as well. Now consider the Hopf ideal $`I/I_1`$ of the quotient Hopf algebra $`GX/I_1.`$ This ideal also has non-zero skew primitive elements (provided $`I_1I`$). Denote by $`I_2/I_1`$ the ideal generated by all skew primitive elements of $`I/I_1,`$ where $`I_2`$ is its preimage with respect to the projection $`GXGX/I_1.`$ Continuing the process we will find a strictly increasing, finite or infinite, chain of Hopf ideals of $`GX:`$
(15)
$$0=I_0I_1I_2\mathrm{}I_n\mathrm{}.$$
###### Definition 3.1.
The length of (15) is called a combinatorial rank of $`H.`$
By definition, the combinatorial rank of any quantum enveloping algebra (with constants) equals one. In the case of zero characteristic the inverse statement is valid as well.
###### Theorem 3.2.
Each character Hopf algebra of the combinatorial rank 1 over a field of zero characteristic is isomorphic to a quantum enveloping algebra with constants of a Lie algebra.
Proof. By definition, $`I`$ is generated by skew primitive elements. These elements as noncommutative polynomials are the quantum operations. Consider one of them, say $`f.`$ Let us decompose $`f`$ into a sum of homogeneous components $`f=f_i.`$ All positive components belongs to $`𝐤X`$ and they are the quantum operations themselves, while the constant component has the form $`\alpha (1g),gG`$ (see \[21, Sec. 3 and Prop. 3.3\]). If $`\alpha 0`$ then we introduce a new quantum variable $`z_f`$ with the parameters $`(id,g).`$ Each $`f_i`$ has a representation through the skew commutator. Indeed, by \[21, Theorem 7.5\] the complete linearization $`f_i^{lin}`$ of $`f_i`$ has the required representation. By the identification of variables in a suitable way in $`f_i^{lin}`$ we get the required representation for $`f_i`$ multiplied by a natural number, $`m_if_i=f_i^{[]}.`$
Now consider a Lie algebra $`𝔤`$ defined by the generators $`x_i,z_f`$ and the relations $`m_i^1f_i^{[]}+z_f=0,`$ with the Lie multiplication in place of the skew commutator. It is clear that $`H`$ is the quantification with constants of $`𝔤.`$
In the same way one may introduce the notion of the combinatorial rank for the braided bigraded Hopf algebras. In this case all braided quantum enveloping algebras are of rank 1, and all braided bigraded algebras of rank 1 are the braided quantification of some Lie algebras.
Now we are ready to define a quantification of arbitrary rank. For this in the definitions of the above section it is necessary to change the requirement that all $`f_i^{}`$ are quantum operations with the following condition.
The set $`F`$ splits in a union $`F=_{j=1}^nF_j`$ such that $`F_1^{}`$ consists of quantum operations; the set $`F_2^{}`$ consists of skew primitive elements of $`GX||F_1^{};`$ the set $`F_3^{}`$ consists of skew primitive elements of $`GX||F_1^{},F_2^{},`$ and so on.
The quantum enveloping algebras of an arbitrary rank are character Hopf algebras also. But it is not clear if any character Hopf algebra is a quantification of some rank of a suitable Lie algebra. It is so if the Hopf algebra is homogeneous and the ground field has a zero characteristic (to appear). Also it is not clear if there exist character Hopf algebras, or braided bigraded Hopf algebras, of infinite combinatorial rank; while it is easy to see that $`I_n=I.`$ Also it is possible to show that $`F_1`$ always contains all relations of a minimal constitution in $`F.`$ For example, each of (11) is of a minimal constitution in (11). Therefore the quantification of arbitrary rank with the identification $`g_i=\text{exp}(h_i)`$ of any (generalized) Kac–Moody algebra $`𝔤,`$ or its nilpotent component $`𝔤^+`$, is always a quantification in the sense of the above section.
4.PBW-generators and crystallisation
The next result yields a PBW basis for the quantum enveloping algebras.
###### Theorem 4.1.
Every character Hopf algebra $`H`$ has a linearly ordered set of constitution homogeneous elements $`U=\{u_i|iI\}`$ such that the set of all products $`gu_1^{n_1}u_2^{n_2}\mathrm{}u_m^{n_m},`$ where $`gG,`$ $`u_1<u_2<\mathrm{}<u_m,`$ $`0n_i<h(i)`$ forms a basis of $`H.`$ Here if $`p_{ii}\stackrel{df}{=}p_{u_iu_i}`$ is not a root of unity then $`h(i)=\mathrm{};`$ if $`p_{ii}=1`$ then either $`h(i)=\mathrm{}`$ or $`h(i)=l`$ is the characteristic of the ground field; if $`p_{ii}`$ is a primitive $`t`$-th root of unity, $`t1,`$ then $`h(i)=t.`$
The set $`U`$ is referred to as a set of PBW-generators of $`H.`$ This theorem easily follows from \[22, Theorem 2\]. Let us recall necessary notions.
Let $`a_1,\mathrm{},a_n`$ be a set of skew primitive generators of $`H,`$ and let $`x_i`$ be the associated quantum variables. Consider the lexicographical ordering of all words in $`x_1>x_2>\mathrm{}>x_n.`$ A non-empty word $`u`$ is called standard if $`vw>wv`$ for each decomposition $`u=vw`$ with non-empty $`v,w.`$ The following properties are well known (see, for example ).
1s. A word $`u`$ is standard if and only if it is greater than each of its ends.
2s. Every standard word starts with a maximal letter that it has.
3s. Each word $`c`$ has a unique representation $`c=u_1^{n_1}u_2^{n_2}\mathrm{}u_k^{n_k},`$ where $`u_1<u_2<\mathrm{}<u_k`$ are standard words (the Lyndon theorem).
4s. If $`u,v`$ are different standard words and $`u^n`$ contains $`v^k`$ as a sub-word, $`u^n=cv^kd,`$ then $`u`$ itself contains $`v^k`$ as a sub-word, $`u=bv^ke.`$
The set of standard nonassociative words is defined as the smallest set $`SL`$ that contains all variables $`x_i`$ and satisfies the following properties.
1) If $`[u]=[[v][w]]SL`$ then $`[v],[w]SL,`$ and $`v>w`$ are standard.
2) If $`[u]=[[[v_1][v_2]][w]]SL`$ then $`v_2w.`$
The following statements are valid as well.
5s. Every standard word has the only alignment of brackets such that the appeared nonassociative word is standard (the Shirshov theorem ).
6s. The factors $`v,w`$ of the nonassociative decomposition $`[u]=[[v][w]]`$ are the standard words such that $`u=vw`$ and $`v`$ has the minimal length ().
###### Definition 4.2.
A super-letter is a polynomial that equals a nonassociative standard word where the brackets mean (2). A super-word is a word in super-letters. By 5s every standard word $`u`$ defines a super-letter $`[u].`$
Let $`D`$ be a linearly ordered Abelian additive group. Suppose that some positive $`D`$-degrees $`d_1,\mathrm{},d_nD`$ are associated to $`x_1,\mathrm{},x_n.`$ We define the degree of a word to be equal to $`m_1d_1+\mathrm{}+m_nd_n`$ where $`(m_1,\mathrm{},m_n)`$ is the constitution of the word. The order and the degree on the super-letters are defined in the following way: $`[u]>[v]u>v;`$ $`\mathrm{D}([u])=\mathrm{D}(u).`$
###### Definition 4.3.
A super-letter $`[u]`$ is called hard in $`H`$ provided that its value in $`H`$ is not a linear combination of values of super-words of the same degree in less than $`[u]`$ super-letters and $`G`$-super-words of a lesser degree.
###### Definition 4.4.
We say that a height of a super-letter $`[u]`$ of degree $`d`$ equals $`h=h([u])`$ if $`h`$ is the smallest number such that: first $`p_{uu}`$ is a primitive $`t`$-th root of unity and either $`h=t`$ or $`h=tl^r,`$ where $`l=`$char(k); and then the value in $`H`$ of $`[u]^h`$ is a linear combination of super-words of degree $`hd`$ in less than $`[u]`$ super-letters and $`G`$-super-words of a lesser degree. If there exists no such number then the height equals infinity.
Clearly, if the algebra $`H`$ is $`D`$-homogeneous then one may omit the underlined parts of the above definitions.
###### Theorem 4.5.
$`([\text{22},\text{Theorem 2}]).`$ The set of all values in $`H`$ of all $`G`$-super-words $`W`$ in the hard super-letters $`[u_i],`$
(16)
$$W=g[u_1]^{n_1}[u_2]^{n_2}\mathrm{}[u_m]^{n_m},$$
where $`gG,`$ $`u_1<u_2<\mathrm{}<u_m,`$ $`n_i<h([u_i])`$ is a basis of $`H.`$
In order to find the set of PBW-generators it is necessary first to include in $`U`$ the values of all hard super-letters, then for each hard super-letter $`[u]`$ of a finite height, $`h([u])=tl^k,`$ to add the values of $`[u]^t,[u]^{tl},\mathrm{}[u]^{tl^{(k1)}},`$ and next for each hard super-letter of infinite height such that $`p_{uu}`$ is a primitive $`t`$-th root of unity to add the value of $`[u]^t.`$
Obviously the set of PBW-generators plays the same role as the basis of the Lie algebra in the PBW theorem does. Nevertheless the k$`[G]`$-bimodule generated by the PBW-generators is not uniquely defined. It depends on the ordering of the main generators, the $`D`$-degree, and under the action of antipode it transforms to a different bimodule of PBW-generators k$`[G]S(U).`$
Another way to construct PBW-generators is connected with the M.Kashiwara crystallisation idea . M.Kashiwara considered the main parameter of the Drinfeld–Jimbo enveloping algebra as the temperature of some physical medium. When the temperature tends to zero the medium crystallises. By this means the ‘crystal’ bases must appear. If we replace $`p_{ij}`$ with zero then $`[u,v]`$ turns into $`uv,`$ while $`[u]`$ turns into $`u.`$
###### Lemma 4.6.
(Bases Crystallisation). Under the above crystallisation the set of PBW-generators constructed in Theorem 4.5 turns into another set of PBW-generators.
Proof. See \[22, Corollary 1\].
###### Lemma 4.7.
(Super-letters Crystallisation). A super-letter $`[u]`$ is hard in $`H`$ if and only if the value of $`u`$ is not a linear combination of lesser words of the same degree and $`G`$-words of a lesser degree.
Proof. See \[22, Corollary 2\].
###### Lemma 4.8.
Let $`B`$ be a set of the super-letters containing $`x_1,\mathrm{},x_n.`$ If each pair $`[u],[v]B,`$ $`u>v`$ satisfies one of the following conditions
$`1)[[u][v]]`$ is not a standard nonassociative word;
$`2)`$ the super-letter $`[[u][v]]`$ is not hard in $`H;`$
$`3)[[u][v]]B,`$
then the set $`B`$ includes all hard in $`H`$ super-letters.
Proof. Let $`[w]`$ be a hard super-letter of minimal degree such that $`[w]B.`$ Then $`[w]=[[u][v]],u>v`$ where $`[u],`$ $`[v]`$ are hard super-letters. Indeed, if $`[u]`$ is not hard then by Lemma 4.7 we have $`u=\alpha _iu_i+S,`$ where $`u_i<u`$ and $`D(u_i)=D(u),`$ $`D(S)<D(u).`$ We have $`uv=\alpha _iu_iv+Sv,`$ where $`u_iv<uv.`$ Therefore by Lemma 4.7, the super-letter $`[w]=[uv]`$ can not be hard in $`H.`$ Contradiction. Similarly, if $`[v]`$ is not hard then $`v=\alpha _iv_i+S,`$ $`v_i<v,`$ $`D(v_i)=D(v),`$ $`D(S)<D(v).`$ Therefore $`uv=\alpha _iuv_i+uS,`$ $`uv_i<uv,`$ and again $`[w]`$ can not be hard.
Thus, according to the choice of $`[w],`$ we get $`[u],[v]B.`$ Since this pair satisfies neither condition 1) nor 2), the condition 3), $`[uv]B,`$ holds. $`\mathrm{}`$
###### Lemma 4.9.
If $`𝐓H`$ is a skew primitive element then
$$𝐓=[u]^h+\alpha _iW_i+\beta _jg_jW_j^{}$$
where $`[u]`$ is a hard super-letter, $`W_i`$ are basis super-words in super-letters less than $`[u],`$ $`D(W_i)=hD([u]),`$ $`D(W_j^{})<hD([u]).`$ Here if $`p_{uu}`$ is not a root of unity then $`h=1;`$ if $`p_{uu}`$ is a primitive $`t`$-th root of unity then $`h=1,`$ or $`h=t,`$ or $`h=tl^k,`$ where $`l`$ is the characteristic.
Proof. Consider an expansion of $`𝐓`$ in terms of the basis (16)
(17)
$$𝐓=\alpha gU+\underset{i=1}{\overset{k}{}}\gamma _ig_iW_i+W^{},\alpha 0,$$
where $`gU,g_iW_i`$ are different basis elements of maximal degree, and $`U`$ is one of the biggest words among $`U,W_i`$ with respect to the lexicographic ordering of words in the super-letters. On basis expansion of tensors, the element $`\mathrm{\Delta }(𝐓)𝐓1g_t𝐓`$ has only one tensor of the form $`gU\mathrm{}`$ and this tensor equals $`gU\alpha (g1).`$ Therefore $`g=1`$ and one may apply \[22, Lemma 13\]. $`\mathrm{}`$
5. Groebner–Shirshov relations systems
Let $`x_1,\mathrm{},x_n`$ be variables that have positive degrees $`d_1,\mathrm{},d_nD.`$ Recall that a Hall ordering of words in $`x_1,\mathrm{},x_n`$ is an order when the words are compared firstly by the degree and then words of the same degree are compared by means of the lexicographic ordering . Consider a set of relations
(18)
$$w_i=f_i,iI,$$
where $`w_i`$ is a word and $`f_i`$ is a linear combination of Hall lesser words. The system (18) is said to be closed under compositions or a Groebner–Shirshov relations system if first none of $`w_i`$ contains $`w_j,`$ $`ijI`$ as a sub-word, and then for each pair of words $`w_k,`$ $`w_j`$ such that some non-empty terminal of $`w_k`$ coincides with an onset of $`w_j,`$ that is $`w_k=w_k^{}v,`$ $`w_j=vw_j^{},`$ the difference (a composition) $`f_kw_j^{}w_k^{}f_j`$ can be reduced to zero in the free algebra by means of a sequence of one sided substitutions $`w_if_i,`$ $`iI.`$
###### Lemma 5.1.
(Diamond Lemma $`[\text{3},\text{5},\text{37}]).`$ If the system $`(\text{18})`$ is closed under compositions then the words that have none of $`w_i`$ as sub-words form a basis of the algebra $`H`$ defined by $`(\text{18}).`$
If none of the words $`w_i`$ has sub-words $`w_j,`$ $`ji,`$ then the converse statement is valid as well. Indeed, any composition by means of substitutions $`w_if_i`$ can be reduced to a linear combination of words that have no sub-words $`w_i.`$ Since $`f_iw_j^{}w_i^{}f_j=`$ $`(f_iw_i)w_j^{}w_i^{}(f_jw_j),`$ this linear combination equals zero in $`H.`$ Therefore all the coefficients have to be zero.
Since Bases Crystallisation Lemma provides the basis that consists of words, the above note gives a way to construct the Groebner–Shirshov relations system for any quantum enveloping algebra.
Let $`H`$ be a character Hopf algebra generated by skew primitive semi-invariants $`a_1,\mathrm{},a_n`$ (or a braided bigraded Hopf algebra generated by graiding homogeneous primitive elements $`a_1,\mathrm{},a_n`$) and let $`x_1,\mathrm{},x_n`$ be the related quantum variables. A non-hard in $`H`$ super-letter $`[w]`$ is referred to as a minimal one if first $`w`$ has no proper standard sub-words that define non-hard super-letters, and then $`w`$ has no sub-words $`u^h,`$ where $`[u]`$ is a hard super-letter of the height $`h.`$
By the Super-letters Crystallisation Lemma, for every minimal non-hard in $`H`$ super-letter $`[w]`$ we may write a relation in $`H`$
(19)
$$w=\alpha _iw_i+\beta _jg_jw_j,$$
where $`w_j,`$ $`w_i<w`$ in the Hall sense, $`D(w_i)=D(w),`$ $`D(w_j)<D(w).`$ In the same way if $`[u]`$ is a hard in $`H`$ super-letter of a finite height $`h`$ then
(20)
$$u^h=\alpha _iu_i+\beta _jg_ju_j,$$
where $`u_j,`$ $`u_i<u^h`$ in the Hall sense, $`D(u_i)=hD(u),`$ $`D(u_j)<hD(u).`$ The relations (14) and the group operation provide the relations
(21)
$$x_ig=\chi ^{x_i}(g)gx_i,g_1g_2=g_3.$$
###### Theorem 5.2.
The set of relations $`(\text{19}),`$ $`(\text{20}),`$ and $`(\text{21})`$ forms a Groebner–Shirshov system that defines $`H.`$ The basis determined by this system in Diamond Lemma coincides with the crystal basis.
Proof. The property 4s implies that none of the left hand sides of (19), (20), (21) contains another one as a sub-word. Therefore by the Bases Crystallisation Lemma it is sufficient to show that the set of all words $`c`$ determined in the Diamond Lemma coincides with the crystal basis. By 3s we have $`c=u_1^{n_1}u_2^{n_2}\mathrm{}u_k^{n_k},`$ where $`u_1<\mathrm{}<u_k`$ is a sequence of standard words. Every word $`u_i`$ define a hard super-letter $`[u_i]`$ since in the opposite case $`u_i,`$ and therefore $`c,`$ contains a sub-word $`w`$ that defines a minimal non-hard super-letter $`[w]`$. In the same way $`n_i`$ does not exceed the height of $`[u_i].`$ $`\mathrm{}`$
###### Lemma 5.3.
In terms of Lemma 4.8 the set of all super-letters $`[[u][v]]`$ that satisfy the condition $`2)`$ contains all minimal non-hard super-letters, but non-hard generators $`x_i.`$
Proof. If $`[w]`$ is a minimal non-hard super-letter then $`[w]=[[u][v]],`$ where $`[u],`$ $`[v]`$ are hard super-letters. By Lemma 4.8 we have $`[u],[v]B,`$ while $`[[u][v]]`$ neither satisfies 1) nor 3). $`\mathrm{}`$
6.Quantification with constants
By means of the Diamond Lemma in some instances the investigation of a quantification with constants can be reduced to one of a simple quantification.
Let $`H_1=x_1,\mathrm{},x_k||F_1`$ be a character Hopf algebra defined by the quantum variables $`x_1,\mathrm{},x_k`$ and the grading homogeneous relations $`\{f=0:fF_1\},`$ while $`H_2=x_{k+1},\mathrm{},x_n||F_2`$ is a character Hopf algebra defined by the quantum variables $`x_{k+1},\mathrm{},x_n`$ and the grading homogeneous relations $`\{h=0:hF_2\}.`$ Consider the algebra $`H=x_1,\mathrm{},x_n||F_1,F_2,F_3,`$ where $`F_3`$ is the following system of relations with constants
(22)
$$[x_i,x_j]=\alpha _{ij}(1g_ig_j),1ik<jn.$$
If the conditions below are met then the character Hopf algebra structure on $`H`$ is uniquely determined:
(23)
$$p_{ij}p_{ji}=1,1ik<jn;\chi ^{x_i}\chi ^{x_j}1\alpha _{ij}=0.$$
Indeed, in this case the difference $`w_{ij}`$ between the left and right hand sides of (22) is a skew primitive semi-invariant of the free enveloping algebra $`Gx_1,\mathrm{},x_n.`$ Consider the ideals of relations $`I_1=`$id$`(F_1)`$ and $`I_2=`$id$`(F_2)`$ of $`H_1`$ and $`H_2`$ respectively. They are, in the present context, Hopf ideals of $`Gx_1,\mathrm{},x_k`$ and $`Gx_{k+1},\mathrm{},x_n,`$ respectively. Therefore $`V=I_1+I_2+`$ $``$k$`w_{ij}`$ is an antipode stable coideal of $`GX.`$ Consequently the ideal generated by $`V`$ is a Hopf ideal. It remains to note that this ideal is generated in $`GX`$ by $`w_{ij}`$ and $`F_1,F_2.`$
###### Lemma 6.1.
Every hard in $`H`$ super-letter belongs to either $`H_1`$ or $`H_2,`$ and it is hard in the related algebra.
Proof. If a standard word contains at least one of the letters $`x_i,`$ $`ik`$ then it has to start with one of them (see s2). If this word contains a letter $`x_j,`$ $`j>k`$ then it has a sub word of the form $`x_ix_j,`$ $`ik<j.`$ Therefore by Lemma 4.7 and relations (22) this word defines a non-hard super-letter. $`\mathrm{}`$
The converse statement is not universally true. In order to formulate the necessary and sufficient conditions let us define partial skew derivatives:
$`(x_j)_i^{}=(x_i)_j^{}=\alpha _{ij}(1g_ig_j),ik<j;`$
$`(vw)_i^{}=(v)_i^{}w+p(x_i,v)v(w)_i^{},ik,v,wGx_{k+1},\mathrm{},x_n;`$
(24) $`(uv)_j^{}=p(v,x_j)(u)_j^{}v+u(v)_j^{},j>k,u,vGx_1,\mathrm{},x_k.`$
###### Lemma 6.2.
All hard in $`H_1`$ or $`H_2`$ super-letters are hard in $`H`$ if and only if $`(h)_i^{}=0`$ in $`H_2`$ for all $`ik,`$ $`hF_2`$ and $`(f)_j^{}=0`$ in $`H_1`$ for all $`j>k,`$ $`fF_1.`$ If these conditions are met then
(25)
$$HH_2_{𝐤[G]}H_1$$
as k$`[G]`$-bimodules and the space generated by the skew primitive elements of $`H`$ equals the sum of these spaces for $`H_1`$ and $`H_2.`$
Proof. By (5) and (24) the following equalities are valid in $`H:`$
(26)
$$0=[x_i,h]=(h)_i^{};0=[f,x_j]=(f)_j^{},ik<j.$$
If all hard in $`H_1`$ or $`H_2`$ super-letters are hard in $`H`$ then $`H_1,`$ $`H_2`$ are sub-algebras of $`H.`$ So (26) proves the necessity of the lemma conditions.
Conversely. Let us consider an algebra $`R`$ defined by the generators $`gG,`$ $`x_1,\mathrm{},x_n`$ and the relations (21), (22). Evidently this system is closed under the compositions. Therefore by Diamond Lemma the set of words $`gvw`$ forms a basis of $`R`$ where $`gG;`$ $`v`$ is a word in $`x_j,`$ $`j>k;`$ and $`w`$ is a word in $`x_i,`$ $`ik.`$ In other words $`R`$ as a bimodule over k$`[G]`$ has a decomposition
(27)
$$R=Gx_{k+1},\mathrm{},x_n_{\text{k}[G]}Gx_1,\mathrm{},x_k.$$
Let us show that the two sided ideal of $`R`$ generated by $`F_2`$ coincides with the right ideal $`I_2R=`$ $`I_2_{\text{k}[G]}Gx_1,\mathrm{},x_k.`$ It will suffice to show that $`I_2R`$ admits left multiplication by $`x_i,`$ $`ik.`$ If $`v`$ is a word in $`x_{k+1},\mathrm{},x_k,`$ $`hF_2,`$ $`rR`$ then $`x_ivhr=[x_i,vh]r+p(x_i,vh)vhx_ir.`$ The second term belongs to $`I_2R,`$ while the first one can be rewritten by (5): $`[x_i,v]h+p(x_i,v)v[x_i,h].`$ Both of these addends belong to $`I_2R`$ since $`[x_i,v]=`$ $`(v)_i^{}`$ $`Gx_{k+1},\mathrm{},x_n`$ and $`[x_i,h]=`$ $`(h)_i^{}I_2.`$
Furthermore, consider a quotient algebra $`R_1=R/I_2R:`$
$$R_1=(Gx_{k+1},\mathrm{},x_n_{\text{k}[G]}Gx_1,\mathrm{},x_k)/(I_2_{\text{k}[G]}Gx_1,\mathrm{},x_k)=$$
$$H_2_{\text{k}[G]}Gx_1,\mathrm{},x_k,$$
where the equality means the natural isomorphism of k$`[G]`$-bimodules.
Along similar lines, the left ideal $`R_1I_1=`$ $`H_2_{\text{k}[G]}I_1`$ of this quotient algebra coincides with the two sided ideal generated by $`F_1.`$ Therefore
$$H=R_1/R_1I_1=H_2_{\text{k}[G]}Gx_1,\mathrm{},x_k/H_2_{\text{k}[G]}I_1=H_2_{\text{k}[G]}H_1.$$
Thus the monotonous restricted $`G`$-words in hard in $`H_1`$ or $`H_2`$ super-letters form a basis of $`H.`$ This, in particular, proves the first statement.
Now let $`T=\alpha _tg_tV_tW_t`$ be the basis decomposition of a skew primitive element, $`g_tG,`$ $`V_tH_2,`$ $`W_tH_1,`$ $`\alpha _t0.`$ We have to show that for each $`t`$ one of the super-words $`V_t`$ or $`W_t`$ is empty. Suppose that it is not so. Among the addends with non-empty $`V_t,`$ $`W_t`$ we choose the largest one in the Hall sense, say $`g_sV_sW_s.`$ Under the basis decomposition of $`\mathrm{\Delta }(T)T1g(T)T`$ the term $`\alpha _sg_sg(V_s)W_sg_sV_s`$ appears and cannot be cancelled with other. Indeed, since the coproduct is homogeneous (see \[22, Lemma 9\]) and since under the basis decomposition the super-words are decreased (see \[22, Lemma 7\]) the product $`\alpha _s(g_sg_s)\mathrm{\Delta }(V_s)\mathrm{\Delta }(W_s)`$ has the only term of the above type. By the same reasons $`\alpha _t(g_tg_t)\mathrm{\Delta }(V_t)\mathrm{\Delta }(W_t)`$ has a term of the above type only if $`V_tV_s`$ and $`W_tW_s`$ with respect to the Hall ordering of the set of all super-words. However, by the choise of $`s,`$ we have $`D(V_sW_s)D(V_tW_t).`$ Hence $`D(V_t)=D(V_s)`$ and $`D(W_t)=D(W_s).`$ In particular $`V_t`$ is not a proper onset of $`V_s.`$ Therefore $`V_t=V_s`$ since otherwise the inequality $`V_t>V_s`$ yields a contradiction $`V_tW_t>V_sW_s`$. The inequality $`W_t>W_s`$ get the same contradiction. Therefore $`V_t=V_s`$ and $`W_t=W_s,`$ in which case $`g_tg(V_t)W_tg_tV_t=`$ $`g_sg(V_s)W_sg_sV_s.`$ Thus $`g_t=g_s`$ and $`t=s.`$ $`\mathrm{}`$
7.Quantification of the classical series
In this section we apply the above general results to the infinite series $`A_n,`$ $`B_n,`$ $`C_n,`$ $`D_n`$ of nilpotent Lie algebras defined by the Serre relations (11). Let $`𝔤`$ be any such Lie algebra.
###### Lemma 7.1.
If a standard word $`u`$ has no sub words of the type
(28)
$$x_i^sx_jx_i^m,\text{where}s+m=1a_{ij}$$
then $`[u]`$ is a hard in $`U_P(𝔤)`$ super-letter.
Proof. Let $`R`$ be defined by the generators $`x_1,\mathrm{},x_n`$ and the relations
(29)
$$x_i^sx_jx_i^m=0,\text{ where }s+m=1a_{ij}.$$
Clearly (29) implies (11) with the skew commutator in place of the Lie operation. Therefore $`R`$ is a homomorphic image of $`U_P(𝔤).`$ The system (29) is closed under compositions since a composition of monomial relations always has the form $`0=0.`$
Let $`u`$ have no sub-words (28). If $`[u]`$ is not hard then, by the Super-letters Crystallisation Lemma, $`u`$ is a linear combination of lesser words in $`U_P(𝔤).`$ Therefore $`u`$ is a linear combination of lesser words in $`R`$ as well. This contradicts the fact that $`u`$ belongs to the Groebner–Shirshov basis of $`R,`$ since every word either belongs to this basis or equals zero in $`R`$. $`\mathrm{}`$
Theorem $`𝐀_𝐧`$. Suppose that $`𝔤`$ is of the type $`A_n,`$ and $`p_{ii}1.`$ Denote by $`B`$ the set of the super-letters given below:
(30)
$$[u_{km}]\stackrel{df}{=}[x_kx_{k+1}\mathrm{}x_m],1kmn.$$
The following statements are valid.
$`1.`$ The values of $`[u_{km}]`$ in $`U_P(𝔤)`$ form a PBW-generators set.
$`2.`$ Each of the super-letters $`(\text{30})`$ has infinite height in $`U_P(𝔤).`$
$`3.`$ The values of all non-hard in $`U_P(𝔤)`$ super-letters equal zero.
$`4.`$ The following relations with $`(\text{21})`$ form the Groebner–Shirshov relations system that determines the crystal basis of $`U_P(𝔤):`$
(31)
$$\begin{array}{cc}[u_0]\stackrel{df}{=}[x_kx_m]=0,& 1k<m1<n;\\ [u_1]\stackrel{df}{=}[x_kx_{k+1}\mathrm{}x_mx_{k+1}]=0,& 1k<mn;\\ [u_2]\stackrel{df}{=}[x_kx_{k+1}\mathrm{}x_mx_kx_{k+1}\mathrm{}x_{m+1}]=0,& 1km<n.\end{array}$$
$`5.`$ If $`p_{11}1`$ then the generators $`x_i,`$ the constants $`1g,`$ $`gG,`$ and, in the case that $`p_{11}`$ is a primitive $`t`$-th root of 1, the elements $`x_i^t,x_i^{tl^k}`$ form a basis of $`𝔤_P=L(U_P(𝔤)).`$ Here $`l`$ is the characteristic of the ground field.
$`6.`$ If $`p_{11}=1`$ then the elements $`(\text{30})`$ and, in the case $`l>0,`$ their $`l^k`$-th powers, together with $`1g,`$ $`gG`$ form a basis of $`𝔤_P.`$
By Corollary 2.4 the relations (11) with a Cartan matrix $`A`$ of type $`A_n`$ admit a quantification if and only if
(32)
$$p_{ii}=p_{11},p_{ii+1}p_{i+1i}=p_{11}^1;p_{ij}p_{ji}=1,ij>1.$$
In this case the quantified relations (11) take up the form
(33) $`x_ix_{i+1}^2=p_{ii+1}(1+p_{i+1i+1})x_{i+1}x_ix_{i+1}p_{ii+1}^2p_{i+1i+1}x_{i+1}^2x_i,`$
(34) $`x_i^2x_{i+1}=p_{ii+1}(1+p_{ii})x_ix_{i+1}x_ip_{ii+1}^2p_{ii}x_{i+1}x_i^2,`$
(35) $`x_ix_j=p_{ij}x_jx_i,ij>1.`$
Let us introduce a congruence $`u_kv`$ on $`GX.`$ This congruence means that the value of $`uv`$ in $`U_P^b(𝔤)`$ belongs to the subspace generated by values of all words with the initial letters $`x_i,ik.`$ Clearly, this congruence admits right multiplication by arbitrary polynomials as well as left multiplication by the independent of $`x_{k1}`$ ones (see (35)). For example, by (33) and (34) we have
(36)
$$x_ix_{i+1}^2_{i+1}0;x_ix_{i+1}x_i_{i+1}\alpha x_i^2x_{i+1},\alpha 0.$$
###### Lemma 7.2.
If $`y=x_i,m+1i>k`$ or $`y=x_i^2,m+1=i>k`$ then
(37)
$$u_{km}y_{k+1}0.$$
Proof. Let $`y=x_{m+1}^2,`$ $`m+1>k.`$ By (36) and (35) we have that $`u_{km}y=`$ $`u_{km1}\underset{¯}{x_mx_{m+1}^2}_{m+1}0.`$ If $`y=x_i`$ and $`m+1i>k`$ then we get $`u_{km}y=`$ $`\alpha u_{ki1}\underset{¯}{x_ix_{i+1}x_i}`$ $`u_{i+2m}_{i+1}\beta \underset{¯}{u_{ki1}x_i^2}u_{i+1m}_{k+1}0`$ by the above case. $`\mathrm{}`$
###### Lemma 7.3.
The brackets in $`[u_{km}]`$ are left-ordered, $`[u_{km}]=[x_k[u_{k+1m}]].`$
Proof. The statement immediately follows from the properties 6s and 2s. $`\mathrm{}`$
###### Lemma 7.4.
If a nonassociative word $`[[u_{km}][u_{rs}]]`$ is standard then $`k=mr;`$ or $`r=k+1,`$ $`ms;`$ or $`r=k,`$ $`m<s.`$
Proof. By definition, $`u_{km}>u_{rs}`$ if and only if either $`k<r;`$ or $`k=r,`$ $`m<s.`$ If $`k=m`$ then $`u_{km}=x_k`$ and $`mr.`$ If $`km`$ then $`[u_{km}]=[x_k[u_{k+1m}]].`$ Therefore $`u_{k+1m}u_{rs},`$ i.e. either $`k+1>r;`$ or $`k+1=r`$ and $`ms.`$ The former case contradicts $`k<r`$ while the latter one does $`k=r.`$ Thus only the possibilities set in the lemma remain. $`\mathrm{}`$
###### Lemma 7.5.
If $`[w]=[[u_{km}][u_{rs}]],`$ $`n1`$ is a standard nonassociative word then the constitution of $`[w]^h`$ does not equal the constitution of any super-word in less than $`[w]`$ super-letters from $`B.`$
Proof. The inequalities at the last column of the following tableaux are valid for all $`[u]B`$ that are less than the super-letters located on the same row, where as above deg$`{}_{i}{}^{}(u)`$ means the degree of $`u`$ in $`x_i.`$
(38)
$$\begin{array}{ccc}[x_ku_{k+1s}]& & \text{deg}_k(u)\text{deg}_{s+1}(u);\\ [x_ku_{rs}],& krk+1& \text{deg}_k(u)\text{deg}_{k+1}(u);\\ [u_{km}u_{k+1s}],& ms& \text{deg}_k(u)\text{deg}_{m+1}(u);\\ [u_{km}u_{ks}],& m<s& \text{deg}_k(u)\text{deg}_{m+1}(u).\end{array}$$
If all super-letters of a super-word $`U`$ satisfy one of these inequalities then $`U`$ does as well. Clearly, no one of the super-letters in the first column satisfies the degree inequality on the same row. Finally, by Lemma 7.4 the first column contains all standard nonassociative words of the type $`[[u_{km}][u_{rs}]].`$ $`\mathrm{}`$
###### Lemma 7.6.
If $`p_{11}1`$ then the values of $`[u_{km}]^h,`$ $`k<m,`$ $`h1`$ are not skew primitive, in particular they are non-zero.
Proof. The sub-algebra generated by $`x_2,\mathrm{}x_n`$ is defined by the Cartan matrix of the type $`A_{n1}.`$ This allows us to use induction on $`n.`$ If $`n=1`$ then the lemma is correct in the sense that $`[u_{km}]^h=x_1^h0.`$
Let $`n>1.`$ If $`k>1`$ then we may use the inductive supposition directly. Consider the decomposition $`\mathrm{\Delta }([u_{1m}])=u^{(1)}u^{(2)}.`$ Since
(39)
$$[u_{1m}]=x_1[u_{2m}]p(x_1,u_{2m})[u_{2m}]x_1,$$
we have
$$\mathrm{\Delta }([u_{1m}])=(x_11+g_1x_1)\mathrm{\Delta }([u_{2m}])$$
(40)
$$p(x_1,u_{2m})\mathrm{\Delta }([u_{2m}])(x_11+g_1x_1).$$
Therefore the sum of all tensors $`u^{(1)}u^{(2)}`$ with deg$`{}_{1}{}^{}(u^{(2)})=1,`$ deg$`{}_{k}{}^{}(u^{(2)})=0,`$ $`k>1`$ has the form $`\epsilon g_1[u_{2m}]x_1,`$ where $`\epsilon =`$ $`1p(x_1,u_{2m})p(u_{2m},x_1)`$ since $`[u_{2m}]g_1=p(u_{2m},x_1)g_1[u_{2m}].`$ By (32) we have $`p_{ij}p_{ji}=1`$ for $`i1>j.`$ Therefore $`\epsilon =1p_{12}p_{21}=`$ $`1p_{11}^10.`$
This implies that in the decomposition $`\mathrm{\Delta }([u_{1m}]^h)=v^{(1)}v^{(2)}`$ the sum of all tensors $`v^{(1)}v^{(2)}`$ with deg$`{}_{1}{}^{}(v^{(2)})=h,`$ deg$`{}_{k}{}^{}(v^{(2)})=0,`$ $`k>1`$ equals $`\epsilon ^h[u_{2m}]^hx_1^h.`$ Thus $`[u_{1m}]^h`$ is not skew primitive in $`U_P(𝔤).`$ $`\mathrm{}`$
Proof of Theorem $`𝐀_𝐧.`$ Let us show firstly that $`B`$ satisfies the conditions of Lemma 4.8. By the Super-letter Crystallisation Lemma $`[w]=[[u_{km}][u_{rs}]]`$ is non-hard if the value of $`u_{km}u_{rs}`$ is a linear combination of lesser words. For $`k=m,`$ $`r=k+1`$ we have $`[w]=[u_{ks}]B.`$ If $`k=m,`$ $`r>k+1`$ then the word $`x_ku_{rs}`$ can be diminished by (34) or (35). If $`km`$ then by Lemma 7.4 the word $`u_{km}u_{rs}`$ has a sub-word of the type $`u_1`$ or $`u_2.`$ Thus we need show only that the values in $`U_P(𝔤)`$ of $`u_1`$ and $`u_2`$ are linear combinations of lesser words.
The word $`u_1`$ has such a representation by Lemma 37. Consider the word $`u_2.`$ Let us show by downward induction on $`k`$ that
(41)
$$u_{km}u_{km+1}_{k+1}\gamma u_{km+1}u_{km},\gamma 0.$$
If $`k=m`$ then one may use (34) with $`i=k.`$ Let $`k<m.`$ Let us transpose the second letter $`x_k`$ of $`u_2`$ as far to the left as possible by (35). We get
$$u_2=\alpha \underset{¯}{x_kx_{k+1}x_k}x_{k+2}\mathrm{}x_mx_{k+1}\mathrm{}x_{m+1},\alpha 0.$$
By (34) we have
$$u_2_{k+1}\beta x_k^2(x_{k+1}x_{k+2}\mathrm{}x_mx_{k+1}\mathrm{}x_{m+1}),\beta 0.$$
Let us apply the inductive supposition to the word in the parentheses. Since $`x_i,`$ $`i>k+1`$ commutes with $`x_k^2`$ according to the formulae (35), we get
$$u_2_{k+1}\gamma \underset{¯}{x_k^2x_{k+1}}x_{k+2}\mathrm{}x_{m+1}x_{k+1}\mathrm{}x_m.$$
Now it remains to replace the underlined sub-word according to (34) and then to transpose the second letter $`x_k`$ to its former position by (35).
Note that for the diminishing of $`u_1,`$ $`u_2`$ we did not use, and we could not use, the relation $`[x_{n1}x_n^2]=0`$ since deg$`{}_{n}{}^{}(u_1)1,`$ deg$`{}_{n}{}^{}(u_2)1.`$
Thus $`B`$ satisfies the conditions of Lemma 4.8. Since none of $`[u_{km}]`$ has sub-words (28), Lemmas 7.1 and 4.8 show that the first statement is correct.
If $`[u_{km}]`$ has a finite height $`h`$ then the value of the polynomial $`[u_{km}]^h`$ in $`U_P(𝔤)`$ is a linear combination of words in hard super-letters that are less than $`[u_{km}].`$ However by Lemma 7.5 this linear combination is trivial, $`[u_{km}]^h=0,`$ since the defining relations are homogeneous. By Lemma 7.6 the second statement is correct for $`p_{11}1`$.
Similarly consider the skew primitive elements. Since both the defining relations and the coproduct are homogeneous, all the homogeneous components of a skew primitive element are skew primitive itself. Therefore it remains to describe all skew primitive elements homogeneous in each $`x_i.`$ Let $`T`$ be such an element. By Lemma 4.9 we have
$$T=[u]^h+\alpha _iW_i,$$
where $`[u]`$ is a hard super-letter, $`u=u_{km},`$ and $`W_i`$ are super-words in less than $`[u]`$ super-letters from $`B`$. By the homogeneity all $`W_i`$ have the same constitution as $`[u_{km}]^h`$ does. However by Lemma 7.5 there exist no such super-words. This means that the only case possible is $`T=[u_{km}]^h.`$ Thus, by Lemma 7.6 the fifth statement is valid as well.
If $`p_{11}=1`$ then $`p_{ij}p_{ji}=p_{ii}=1`$ for all $`i,j.`$ So we are under the conditions of Example 1, that is $`U_P^b(𝔤)`$ is the universal enveloping algebra of the colour Lie algebra $`𝔤^{col}.`$ Further, $`[u_{km}]𝔤^{col}`$ and $`[u_{km}]`$ are linearly independent in $`𝔤^{col}`$ since they are hard super-letters and no one of them can be a linear combination of the lesser ones. Let us complete $`B`$ to a homogeneous basis $`B^{}`$ of $`𝔤^{col}.`$ Then by the PBW theorem for the colour Lie algebras the products $`b_1^{n_1}\mathrm{}b_k^{n_k},`$ $`b_1<\mathrm{}<b_k`$ form a basis of $`U(𝔤^{col})=`$ $`U_P^b(𝔤).`$ However, the monotonous restricted words in $`B`$ form a basis of $`U_P^b(𝔤)`$ also. Thus $`B^{}=B`$ and all hard super-letters have the infinite height.
In particular, we get that the second statement is valid in complete extent. Moreover, if $`p_{11}=1`$ then $`p(u_{km},u_{km})=1,`$ thus for $`l=0`$ all homogeneous skew primitive elements became exhausted by $`[u_{km}],`$ while for $`l>0`$ the powers $`[u_{km}]^{l^k}`$ are added to them (of course, here $`l2`$ since $`1p_{ii}=1`$).
So we have proved all statements, but the third and fourth ones. These statements will follow Theorem 5.2 and Lemma 5.3 if we prove that all non-hard super-letters $`[[u_{km}][u_{rs}]]`$ equal zero in $`U_P(𝔤).`$ By the homogeneous definition, $`[[u_{km}][u_{rs}]]`$ is a linear combination of super-words in lesser hard super-letters. However, by Lemma 7.5, there exist no such super-words of the same constitution. Therefore, by the homogeneity, the above linear combination equals zero. $`\mathrm{}`$
Theorem $`𝐁_𝐧`$. Let $`𝔤`$ be of the type $`B_n,`$ and $`p_{ii}1,`$ $`1i<n,`$ $`p_{nn}^{[3]}0.`$ Denote by $`B`$ the set of the super-letters given below:
(42)
$$\begin{array}{cc}[u_{km}]\stackrel{df}{=}[x_kx_{k+1}\mathrm{}x_m],& 1kmn;\\ [w_{km}]\stackrel{df}{=}[x_kx_{k+1}\mathrm{}x_nx_n\mathrm{}x_m],& 1k<mn.\end{array}$$
The following statements are valid.
$`1.`$ The values of $`(\text{42})`$ in $`U_P(𝔤)`$ form the PBW-generators set.
$`2.`$ Every super-letter $`[u]B`$ has infinite height in $`U_P(𝔤).`$
$`3.`$ The relations $`(\text{21})`$ with the following ones form a Groebner–Shirshov system that determines the crystal basis of $`U_P(𝔤).`$
(43)
$$\begin{array}{cc}[u_0]\stackrel{df}{=}[x_kx_m]=0,& 1k<m1<n;\\ [u_1]\stackrel{df}{=}[u_{km}x_{k+1}]=0,& 1k<mn,kn1;\\ [u_2]\stackrel{df}{=}[u_{km}u_{km+1}]=0,& 1km<n;\\ [u_3]\stackrel{df}{=}[w_{km}x_{k+1}]=0,& 1k<mn,km2;\\ [u_4]\stackrel{df}{=}[w_{kk+1}x_{k+2}]=0,& 1k<n1;\\ [u_5]\stackrel{df}{=}[w_{km}w_{km1}]=0,& 1k<m1n1;\\ [u_6]\stackrel{df}{=}[u_{kn}^2x_n]=0,& 1k<n.\end{array}$$
$`4.`$ If $`p_{11}1`$ then the generators $`x_i`$ and their powers $`x_i^t,x_i^{tl^k},`$ such that $`p_{ii}`$ is a primitive $`t`$-th root of 1, together with the constants $`1g,`$ $`gG`$ form a basis of $`𝔤_P=L(U_P(𝔤)).`$ Here $`l`$ is the characteristic of the ground field.
$`5.`$ If $`p_{nn}=p_{11}=1`$ then the elements $`(\text{42})`$ and, for $`l>0,`$ their $`l^k`$-th powers, together with $`1g,`$ $`gG`$ form a basis of $`𝔤_P.`$ If $`p_{nn}=p_{11}=1`$ then $`[u_{kn}]^2,`$ $`[u_{kn}]^{2l^k}`$ are added to them.
Recall that in the case $`B_n`$ the algebra $`U_P^b(𝔤)`$ is defined by (33), (34), (35) where in (33) the last relation, $`i=n1,`$ is replaced with
(44)
$$x_{n1}x_n^3=p_{n1n}p_{nn}^{[3]}x_nx_{n1}x_n^2p_{n1n}^2p_{nn}p_{nn}^{[3]}x_n^2x_{n1}x_n+p_{n1n}^3p_{nn}^3x_n^3x_{n1}.$$
By Corollary 2.4 we get the existence conditions
(45)
$$p_{ii}=p_{11},p_{ii+1}p_{i+1i}=p_{11}^1=p_{nn}^2,1in1;p_{ij}p_{ji}=1,ij>1.$$
The relations (33) and (44) show that
(46)
$$x_ix_{i+1}^2_{i+1}0,i<n1;x_{n1}x_n^3_n0,$$
while the relations (34) imply
(47)
$$x_ix_{i+1}x_i_{i+1}\alpha x_i^2x_{i+1},\alpha 0.$$
By means of these relations and (35), (44) we have
(48)
$$x_{n2}x_{n1}\underset{¯}{x_n^2x_{n1}x_n}_{n1}0.$$
###### Lemma 7.7.
The brackets in $`[w_{km}]`$ are set by the recurrence formulae:
(49)
$$\begin{array}{cc}[w_{km}]=[x_k[w_{k+1m}]],& \text{if}1k<m1<n;\\ [w_{kk+1}]=[[w_{kk+2}]x_{k+1}],& \text{if}1k<n.\end{array}$$
Here by the definition $`w_{kn+1}=u_{kn}.`$
Proof. It is enough to use the property 6s and then 1s and 2s. $`\mathrm{}`$
###### Lemma 7.8.
The nonassociative word $`[[w_{km}][w_{rs}]]`$ is standard only in the following two cases: $`1)sm>k+1=r;`$ $`2)s<m,`$ $`r=k.`$
Proof. If $`[[w_{km}][w_{rs}]]`$ is standard then $`w_{km}>w_{rs}`$ and by (49) either $`w_{k+1}w_{rs},`$ or $`m=k+1`$ and $`x_{k+1}w_{rs}.`$ The inequality $`w_{km}>w_{rs}`$ is correct only in two cases: $`k<r`$ or $`k=r,m>s.`$ We get four possibilities: $`1)k<r,`$ $`k<m1,`$ $`w_{k+1m}w_{rs};`$ $`2)k<r,`$ $`m=k+1,`$ $`x_{k+1}w_{rs};`$ $`3)k=r,`$ $`m>s,`$ $`k<m1,`$ $`w_{k+1m}w_{rs};`$ $`4)k=r,`$ $`m>s,`$ $`m=k+1,`$ $`x_{k+1}w_{rs}.`$ Only the first and third ones are consistent since in the second case $`x_{k+1}w_{rs}`$ implies $`k+1>r,`$ while in the fourth case $`r<s`$ and $`k=r<s<m=k+1.`$ If now we decode $`w_{k+1m}w_{rs}`$ in the first and third cases, we get the two possibilities mentioned in the lemma. $`\mathrm{}`$
###### Lemma 7.9.
The nonassociative word $`[[u_{km}][w_{rs}]]`$ is standard only in the following two cases: $`1)k=r;`$ $`2)k=m<r.`$
Proof. The inequality $`u_{km}>w_{rs}`$ means $`kr.`$ Since $`[u_{km}]=[x_k[u_{k+1m}]],`$ for $`km`$ we get $`u_{k+1m}w_{rs},`$ so $`k+1>r`$ and $`k=r.`$ If $`k=mr`$ then $`x_m>w_{rs}`$ and $`m<r.`$ $`\mathrm{}`$
###### Lemma 7.10.
The nonassociative word $`[[w_{km}][u_{rs}]]`$ is standard only in the following two cases: $`1)r=k+1<m;`$ $`2)r=k+1=m=s.`$
Proof. The inequality $`w_{km}>u_{rs}`$ implies $`r>k.`$ If $`k<m1`$ then by the first formula (49) we have $`w_{k+1m}u_{rs}`$ that is equivalent to $`k+1r.`$ Therefore $`r=k+1<m.`$ If $`k=m1`$ then by the second formula (49) we get $`x_{k+1}u_{rs},`$ i.e. either $`k+1>r`$ or $`k+1=r=s.`$ The former case contradicts $`r>k`$ while the latter one is mentioned in the lemma. $`\mathrm{}`$
###### Lemma 7.11.
If $`[u],[v]B`$ then one of the statements below is correct.
$`1)[[u][v]]`$ is not a standard nonassociative word;
$`2)uv`$ contains a sub-word of one of the types $`u_0,u_1,u_2,u_3,u_4,u_5,u_6;`$
$`3)[[u][v]]B.`$
Proof. The proof results from Lemmas 7.4, 7.8, 7.9, 7.10. $`\mathrm{}`$
###### Lemma 7.12.
If a super-word $`W`$ equals one of the super-letters $`[u_1]`$$`[u_6]`$ or $`[u_{km}]^h,`$ $`[w_{km}]^h,`$ $`h1`$ then its constitution does not equal the constitution of any super-word in less than $`W`$ super-letters from $`B.`$
Proof. The proof is akin to Lemma 7.5 with the following tableaux:
(50)
$$\begin{array}{cc}[u_{km}],[u_{km}x_{k+1}],[u_{km}u_{km+1}]& \text{deg}_k(u)\text{deg}_{m+1}(u);\\ [w_{km}],[w_{km}x_{k+1}],[w_{km}w_{km1}]& 2\text{deg}_k(u)\text{deg}_{m1}(u);\\ [w_{kk+1}x_{k+2}]& \text{deg}_k(u)=0;\\ [u_{kn}^2x_n]& \text{deg}_k(u)\text{deg}_n(u).\end{array}$$
$`\mathrm{}`$
###### Lemma 7.13.
If $`y=x_i,m1i>k`$ or $`y=x_i^2,m1=i>k`$ then
(51)
$$w_{km}y_{k+1}0.$$
Proof. If $`i<m1`$ then by means of (35) it is possible to permute $`y`$ to the left beyond $`x_n^2`$ and use Lemma 37 with $`m^{}=n1.`$ If $`y=x_i^2,`$ $`m1=i>k`$ then by the above case, $`i<m1,`$ we get
(52)
$$w_{km}y=w_{km+1}\underset{¯}{x_mx_{m1}^2}=\underset{¯}{w_{km+1}x_{m1}}(\alpha x_mx_{m1}+\beta x_{m1}x_m)_{k+1}0,$$
where for $`m=n`$ by definition $`w_{kn+1}=`$ $`u_{kn},`$ and $`u_{kn}x_{n1}_{n1}0.`$
If $`y=x_i,i=m>k`$ then for $`m=n`$ one may use the second equality (46). For $`m<n`$ we have $`w_{km}y=w_{km+1}y_1`$ where $`y_1=x_m^2.`$ Therefore for $`k<n1`$ we may use (52) with $`m+1`$ in place of $`m.`$ For $`k=n1`$ we have $`w_{km}x_n=x_{n1}x_n^3_n0.`$
Finally, if $`y=x_i,i>m>k`$ then by (35) we have $`w_{km}y=\alpha w_{ki+1}\underset{¯}{x_ix_{i1}x_i}v.`$ For $`i=n`$ one may use (48), while for $`i<n,`$ changing the underlined word according to (33), we may use the above considered cases: $`m^{}1=i^{},`$ where $`m^{}=i+1,`$ $`i^{}=i;`$ and $`i^{}<m^{}1,`$ where $`m^{}=i+1,`$ $`i^{}=i1.`$ $`\mathrm{}`$
Another interesting relation appears if we multiply (44) by $`x_{n1}`$ from the left and subtract (34) with $`i=n1`$ multiplied from the right by $`x_n^2:`$
(53)
$$x_{n1}x_nx_{n1}x_n^2_n\alpha x_{n1}x_n^2x_{n1}x_n,$$
in which case $`\alpha =p_{n1n}p_{nn}^{[3]}0.`$
###### Lemma 7.14.
For $`k<s<mn`$ the following relation is valid.
(54)
$$w_{km}w_{ks}_{k+1}\epsilon w_{ks}w_{km},\epsilon 0.$$
Proof. Let us use downward induction on $`k.`$ For this we first transpose the second letter $`x_k`$ of $`w_{km}w_{ks}`$ as far to the left as possible by means of (35), and then change the onset $`x_kx_{k+1}x_k`$ according to (47). We get
(55)
$$w_{km}w_{ks}_{k+1}\alpha x_k^2(w_{k+1m}w_{k+1s}),\alpha 0.$$
For $`k+1<s`$ we apply the inductive supposition to the word in the parentheses and then by (47) and (35) transpose $`x_k`$ to its former position.
The case $`k+1=s,`$ the basis of the induction on $`k,`$ we prove by downward induction on $`s.`$
Let $`k+1=s=n1.`$ Then $`m=n.`$ Let us show firstly that
(56)
$$\underset{¯}{x_{n1}x_n^2x_{n1}x_n}x_nx_{n1}_n\alpha x_{n1}x_n^2x_{n1}^2x_n^2+\beta x_{n1}x_nx_{n1}^2x_n^3,\alpha 0.$$
For this in the left hand side we transpose the first letter $`x_n`$ by means of (53) to the penultimate position, and then replace the ending $`x_n^3x_{n1}`$ by (44). We get a linear combination of three words. One of them equals the second word of (56), while two other have the following forms.
$$x_{n1}x_n\underset{¯}{x_{n1}x_nx_{n1}}x_n^2,\underset{¯}{x_{n1}x_nx_{n1}x_n^2}x_{n1}x_n.$$
The former word by (34) transforms into the form (56). The latter one, after the application of (53) and the replacing of $`x_{n1}x_nx_{n1}`$ by (34), will have an additional term $`\underset{¯}{x_{n1}x_n^3}x_{n1}^2x_n`$ to which it is possible to apply (46). The direct calculation of the coefficients shows that $`\alpha =p_{n1n}p_{nn}0.`$
Now let us multiply (56) by $`x_{n2}^2`$ from the left and use (34) with $`i=n2.`$ We get that $`w_{n2n}w_{n2n1}`$ with respect to $`_{n1}`$ equals
(57)
$$\gamma x_{n2}x_{n1}x_n^2\underset{¯}{x_{n2}x_{n1}^2}x_n^2+\delta x_{n2}x_{n1}x_n\underset{¯}{x_{n2}x_{n1}^2}x_n^3,\gamma 0.$$
Let us apply (46) and then (47) and (46) to the second word. We get that this word with respect to $`_{n1}`$ equals zero. The first word after application of (34) takes up the form
$$\epsilon w_{n2n1}w_{n2n}+\epsilon ^{}\underset{¯}{w_{n2n}x_{n1}^2}x_{n2}x_n^2,\epsilon 0.$$
Thus, by Lemma 51, the basis of the induction on $`s`$ is proved.
Let us carry out the inductive step. Let $`k+1=s<n1.`$ If $`m>s+1=k+2`$ then by the inductive supposition on $`s`$ we may write
$$w_{km}w_{ks}=(w_{km}w_{kk+2})x_{k+1}_{k+1}\alpha w_{kk+2}w_{km}x_{k+1}=$$
(58)
$$\beta w_{kk+2}\underset{¯}{x_kx_{k+1}x_{k+2}x_{k+1}}w_{k+3m}.$$
Taking into account (51) we may neglect the words starting with $`x_{k+1}^2,`$ $`x_{k+2}`$ while transforming the underlined part:
(59)
$$x_k\underset{¯}{x_{k+1}x_{k+2}x_{k+1}}\gamma \underset{¯}{x_kx_{k+1}^2}x_{k+2}\delta x_{k+1}x_kx_{k+1}x_{k+2}.$$
In this way (58) is transformed into (54).
If $`m=s+1=k+2<n`$ then the relation (55) takes up the form
$$w_{km}w_{ks}_{k+1}\alpha x_k^2(w_{k+1k+2}w_{k+1k+3})x_{k+2}x_{k+1}.$$
Let us apply the inductive supposition with $`k^{}=k+1,`$ $`s^{}=k+2,`$ $`m^{}=k+3`$ to the word in the parentheses. We get
$$w_{km}w_{ks}_{k+1}\alpha \epsilon ^1x_k^2w_{k+1k+3}w_{k+1k+3}\underset{¯}{x_{k+2}^2x_{k+1}},$$
or after an evident replacement
$$w_{km}w_{ks}_{k+1}\gamma x_k^2w_{k+1k+3}w_{k+1k+2}x_{k+1}x_{k+2}+\delta x_k^2w_{k+1k+3}^2x_{k+1}x_{k+2}^2.$$
In both terms we may transpose one letter $`x_k`$ to its former position by means of (47) and (35). We get
(60)
$$w_{km}w_{ks}_{k+1}\gamma ^{}\underset{¯}{w_{kk+3}w_{kk+1}}x_{k+2}+\delta ^{}w_{kk+3}^2x_{k+1}x_{k+2}^2.$$
It is possible to apply (54) with $`m^{}=k+3,`$ $`s^{}=k+1`$ to the first term since the case $`m>s+1`$ is completely considered. Therefore it is enough to show that the second term equals zero with respect to $`_{k+1}.`$ When we transpose the third letter $`x_{k+1}`$ as far to the left as possible we get the word
(61)
$$w_{kk+3}\underset{¯}{x_kx_{k+1}x_{k+2}x_{k+1}}w_{k+3k+3}x_{k+2}^2.$$
Taking into account (51) we may neglect the words starting with $`x_{k+1}`$ while transforming the underlined part:
(62)
$$x_k\underset{¯}{x_{k+1}x_{k+2}x_{k+1}}x_{k+2}\underset{¯}{x_kx_{k+1}^2}x_{k+2}x_{k+1}x_kx_{k+1}.$$
Therefore the word (61) equals $`w_{kk+1}\underset{¯}{w_{kk+3}x_{k+2}^2}`$ with respect to $`_{k+1}`$ and it remains only to apply Lemma 51 twice. $`\mathrm{}`$
###### Lemma 7.15.
The set $`B`$ satisfies the conditions of Lemma $`\text{4.8}.`$
Proof. By Lemmas 7.11 and 4.7 it is sufficient to show that in $`U_P^b(𝔤)`$ all words of the form $`u_0,\mathrm{},u_6`$ are linear combinations of lesser ones. The words $`u_0`$ are diminished by (35). The words $`u_1,u_2`$ have been presented in this way, without using $`[x_{n1}x_n^2]=0,`$ in the proof of the above theorem. The relation (51) shows that $`u_3_{k+1}0,`$ $`u_4_{k+1}0.`$ Lemma 54 with $`s=m1`$ yields the necessary representation for $`u_5.`$
Let us prove by downward induction on $`k`$ that
$$u_6\stackrel{df}{=}u_{kn}^2x_n_{k+1}\epsilon u_{kn}x_nu_{kn},\epsilon 0.$$
For $`k=n1`$ this equality takes up the form (53). Let $`k<n1.`$ Let us transpose the second letter $`x_k`$ of $`u_{kn}^2x_n`$ as far to the left as possible by means of (35) and then apply (33). We get
$$u_{kn}^2x_n_{k+1}\alpha x_k^2(u_{k+1n}^2x_n),\alpha 0.$$
We may apply the inductive supposition to the term in the parentheses and then by (33), (35) transpose one of $`x_k`$’s to its former position. $`\mathrm{}`$
###### Lemma 7.16.
If $`p_{11}1`$ then the values of polynomials $`[v]^h,`$ where $`[v]B,`$ $`vx_i`$ $`h1`$ are not skew primitive, in particular, they are non-zero.
Proof. Note that for $`n>2`$ the sub-algebra generated by $`x_2,\mathrm{}x_n`$ is defined by the Cartan matrix of the type $`B_{n1}.`$ This allows us to carry out the induction on $`n`$ with additional supposition that the statements 1 and 2 of Theorem $`B_n`$ are valid for lesser values of $`n.`$ It is convenient formally consider the sub-algebras $`x_i`$ as algebras of the type $`B_1.`$ In this case for $`n=1`$ the lemma and the statements 1 and 2 are correct in the evident way. If $`v`$ starts with $`x_kx_1`$ then we may directly use the inductive supposition. If $`v=u_{1m},`$ one may literally repeat the arguments of Lemma 7.6 starting at the formula (39). Let $`v=w_{1m}.`$ If $`m>2`$ then by Lemma 7.7 we have $`w_{1m}=[x_1[w_{2m}]].`$ This provides a possibility to repeat the same arguments of Lemma 7.6 with $`w`$ in place of $`u.`$
Consider the last case $`v=w_{12}.`$ By Lemma 7.7 we have
(63)
$$[w_{12}]=[w_{13}]x_2p(w_{13},x_2)x_2[w_{13}],$$
(64)
$$[w_{13}]=x_1[w_{23}]p(x_1,w_{23})[w_{23}]x_1.$$
Applying the coproduct first to (64) then to (63) we may find the sum $`\mathrm{\Sigma }`$ of all tensors $`w^{(1)}w^{(2)}`$ of $`\mathrm{\Delta }([w_{12}])`$ with deg$`{}_{1}{}^{}(w^{(2)})=1,`$ deg$`{}_{k}{}^{}(w^{(2)})=0,`$ $`k>1`$ (in much the same way as (40)):
$`\mathrm{\Sigma }=(\epsilon g_1[w_{23}]x_1)(x_21)p(w_{13},x_2)(x_21)(\epsilon g_1[w_{23}]x_1)=`$
(65) $`\epsilon g_1([w_{23}]x_2p(w_{13},x_2)p(x_2,x_1)x_2[w_{23}])x_1.`$
For $`n>2,`$ taking into account first the bicharacter property of $`p,`$ then the equality $`[x_2[w_{23}]]=`$ $`x_2[w_{23}]`$ $`p(x_2,w_{23})[w_{23}]x_2,`$ and next the following relations $`p_{ij}p_{ji}=1,`$ $`ij>1;`$ $`p_{11}^1=`$ $`p_{12}p_{21}=`$ $`p_{22}^1=`$ $`p_{23}p_{32},`$ we may write
(66)
$$\mathrm{\Sigma }=\epsilon g_1(p(w_{13},x_2)p_{21}[x_2w_{23}]+(1p_{11}^1)[w_{23}]x_2)x_1.$$
Consider the left hand side of this tensor on applying the inductive supposition. Note that $`x_2w_{23}`$ is a standard word and $`[x_2w_{23}]`$ equals $`[x_2[w_{23}]].`$ This super-letter is non-hard in $`U_P(𝔤)`$ since $`x_2w_{23}`$ contains the sub-word $`x_2^2x_3.`$ Thus $`[x_2w_{23}]`$ is a linear combination of monotonous non-decreasing super-words in lesser super-letters. Among these super-words there is no $`[w_{23}]x_2`$ since $`x_2>x_2w_{23}.`$ On the other hand, $`[w_{23}]x_2`$ is a monotonous non-decreasing super-word and hence its value in $`U_P(𝔤)`$ is a basis element. Therefore for $`n>2`$ the left hand side $`W`$ of $`\mathrm{\Sigma }`$ is non-zero.
For $`n=2,`$ by the definition $`w_{23}=x_2,`$ $`w_{13}=x_1x_2,`$ and the equality (65) takes up the form $`\mathrm{\Sigma }=\epsilon g_1(1p_{12}p_{22}p_{21})x_2^2x_1.`$ Since $`1p_{11}^1=p_{12}p_{21}=p_{22}^2,`$ we get $`(1p_{12}p_{22}p_{21})=1p_{22}^10.`$ Therefore in this case $`\mathrm{\Sigma }0`$ as well.
By \[22, Corollary 10\] the sub-algebra generated by $`x_2,\mathrm{},x_n`$ has no zero divisors. In particular $`W^h0`$ and $`\mathrm{\Sigma }^h0`$ in any case.
It remains to note that for $`n>1`$ the sum of all tensors $`w^{(1)}w^{(2)}`$ of $`\mathrm{\Delta }([w_{12}]^h)`$ such that deg$`{}_{1}{}^{}(w^{(2)})=h,`$ deg$`{}_{k}{}^{}(w^{(2)})=0,`$ $`k>1`$ equals $`\mathrm{\Sigma }^h,`$ hence $`[w_{12}]^h`$ can not be skew-primitive. $`\mathrm{}`$
Proof of Theorem $`B_n.`$ Since none of $`u_{km},`$ $`w_{km}`$ contains sub-words (28), Lemmas 7.15, 7.1, 4.8 imply the first statement.
If $`[v]B`$ is of finite height then by Lemma 7.12 and the homogeneous version of Definition 4.4 we have $`[v]^h=0.`$ For $`p_{11}1`$ this contradicts Lemma 7.16.
Along similar lines, by Lemma 4.9, every skew primitive homogeneous element has the form $`[v]^h.`$ This, together with Lemma 7.16, proves the fourth statement and, for $`p_{11}1,`$ the second one too.
If $`p_{11}=1`$ then by (45) we have $`p_{nn}^2=1,`$ $`p_{ii}=1,`$ $`i<n.`$ Besides, $`p_{ij}p_{ji}=1`$ for all $`i,j.`$ This means that the skew commutator is a quantum operation. Hence all elements of $`B`$ are skew primitive. In the case $`p_{nn}=1`$ these elements span a colour Lie algebra, while in the case $`p_{nn}=1`$ they span a colour Lie super-algebra. Now as in Theorem $`A_n,`$ we may use the $`PBW`$-theorem for the colour Lie super-algebras.
The third statement will follow Theorem 5.2 and Lemmas 5.3, 7.11 if we prove that all super-letters (43) are zero in $`U_P(𝔤).`$ We have already proved that these super-letters are non-hard. Therefore it remains to use the homogeneous version of Definition 4.3 and Lemma 7.12. $`\mathrm{}`$
Theorem $`𝐂_𝐧`$. Suppose that $`𝔤`$ is of the type $`C_n,`$ and $`p_{ii}1,`$ $`1in,`$ $`p_{n1n1}^{[3]}0.`$ Denote by $`B`$ the set of the following super-letters:
(67)
$$\begin{array}{cccc}[u_{km}]& \stackrel{df}{=}& [x_kx_{k+1}\mathrm{}x_m],& 1kmn;\\ [v_{km}]& \stackrel{df}{=}& [x_kx_{k+1}\mathrm{}x_nx_{n1}\mathrm{}x_m],& 1k<m<n;\\ [v_k]& \stackrel{df}{=}& [u_{kn1}u_{kn}],& 1k<n.\end{array}$$
The statements given below are valid.
$`1.`$ The values of the super-letters $`(\text{67})`$ in $`U_P(𝔤)`$ form the PBW-generators set.
$`2.`$ Each of these super-letters has the infinite height in $`U_P(𝔤).`$
$`3.`$ The following relations with $`(\text{21})`$ form a Groebner–Shirshov system that determines the crystal basis of $`U_P(𝔤).`$
(68)
$$\begin{array}{cccc}[u_0]& \stackrel{df}{=}& [x_kx_m]=0,& 1k<m1<n;\\ [u_1]& \stackrel{df}{=}& [u_{km}x_{k+1}]=0,& 1k<mn,(k,m)(n2,n);\\ [u_2]& \stackrel{df}{=}& [u_{km}u_{km+1}]=0,& 1km<n1;\\ [w_3]& \stackrel{df}{=}& [v_{km}x_{k+1}]=0,& 1k<m<n,km2;\\ [w_4]& \stackrel{df}{=}& [v_{kk+1}x_{k+2}]=0,& 1k<n1;\\ [w_5]& \stackrel{df}{=}& [v_{km}v_{km1}]=0,& 1k<m1n1;\\ [w_6]& \stackrel{df}{=}& [u_{kn1}^3x_n]=0,& 1k<n.\end{array}$$
$`4.`$ If $`p_{11}1`$ then the generators $`x_i`$ and their powers $`x_i^t,x_i^{tl^k},`$ such that $`p_{ii}`$ is a primitive $`t`$-th root of 1 together with the constants $`1g,`$ $`gG`$ form a basis of $`𝔤_P=L(U_P(𝔤)).`$ Here $`l`$ is the characteristic of the ground field.
$`5.`$ If $`p_{11}=1`$ then the elements $`(\text{67})`$ and in the case of prime characteristic $`l`$ theirs $`l^k`$-th powers, together with the constants $`1g,`$ $`gG`$ form a basis of $`𝔤_P.`$
In the case $`C_n`$ the algebra $`U_P^b(𝔤)`$ is defined by the same relations (33), (34), (35), where in (34) the last relation, $`i=n1,`$ is replaced with
(69) $`x_{n1}^3x_n=p_{n1n}p_{n1n1}^{[3]}x_{n1}^2x_nx_{n1}+`$
$`p_{n1n}^2p_{n1n1}p_{n1n1}^{[3]}x_{n1}x_nx_{n1}^2+p_{n1n}^3p_{n1n1}^3x_nx_{n1}^3.`$
By Corollary 2.4 we get the existence conditions
$`p_{ii}=p_{11},p_{i1i}p_{ii1}=p_{11}^1,1<i<n,`$
(70) $`p_{n1n}p_{nn1}=p_{nn}^1=p_{n1n1}^2;p_{ij}p_{ji}=1,ij>1.`$
Therefore the following relations are correct
(71) $`x_ix_{i+1}^2_{i+1}0,1i<n;`$
(72) $`x_ix_{i+1}x_i_{i+1}\alpha x_i^2x_{i+1},1i<n1,\alpha 0;`$
(73) $`x_{n1}x_nx_{n1}^2_n\alpha x_{n1}^3x_n+\beta x_{n1}^2x_nx_{n1},\alpha ,\beta 0.`$
The left multiplication by $`x_{n2}`$ of the last relation implies
(74)
$$x_{n2}x_{n1}x_nx_{n1}^2_{n1}0.$$
###### Lemma 7.17.
The brackets in $`[v_{km}],[v_k]`$ are set according to the following recurrence formulae, where by the definition $`v_{kn}=u_{kn}.`$
(75)
$$\begin{array}{cccc}[v_{km}]& =& [x_k[v_{k+1m}]],& \text{if}1k<m1<n1;\\ [v_{kk+1}]& =& [[v_{kk+2}]x_{k+1}],& \text{if}1k<n1;\\ [v_k]& =& [[u_{kn1}][u_{kn}]],& \text{if}1k<n.\end{array}$$
Proof. It is enough to use the properties 6s, 1s and 2s. $`\mathrm{}`$
###### Lemma 7.18.
If $`[u],[v]B`$ then one of the following statements is valid.
$`1)[[u][v]]`$ is not a standard nonassociative word;
$`2)uv`$ contains a sub-word of one of the types $`u_0,u_1,u_2,w_3,w_4,w_5,w_6;`$
$`3)[[u][v]]B.`$
Proof. The first two formulae (75) coincide with (49) up to replacement of $`v`$ with $`w`$ provided $`k+1n>m.`$ Obviously for $`m<n`$ the inequality $`v_{km}>v_{rs}`$ is equivalent to $`w_{km}>w_{rs},`$ while $`v_{km}>u_{rs}`$ is equivalent to $`w_{km}>w_{rs}.`$ Hence Lemmas (7.8), (7.9), (7.10) are still valid under the replacement of $`w`$ with $`v:`$
(76)
$$\begin{array}{ccc}[[v_{km}][v_{rs}]]& \text{ is standard }& sm>k+1=r(s<m\&r=k);\\ [[u_{km}][v_{rs}]]& \text{ is standard }& k=rk=m<r;\\ [[v_{km}][u_{rs}]]& \text{ is standard }& r=k+1<mr=k+1=m=s.\end{array}$$
Further, $`v_k>v_r`$ if and only if $`k<r,`$ and under this condition $`[[v_k][v_r]]`$ is not standard since $`u_{kn}>u_{rn1}u_{rn}.`$
In a similar manner $`v_k>u_{rm}`$ is equivalent to $`k<r,`$ while $`v_k>v_{rm}`$ is equivalent to $`kr.`$ Therefore none of the words $`[[v_k][u_{rm}]],`$ $`[[v_k][v_{rm}]]`$ is standard since $`u_{kn}>u_{rm}`$ and $`u_{kn}>v_{rm},`$ respectively.
For the remaining two cases we have only two possibilities
(77)
$$\begin{array}{ccc}[[u_{km}][v_r]]& \text{ is standard }& r=km<n;\\ [[v_{km}][v_r]]& \text{ is standard }& r=k+1\&k<m1.\end{array}$$
The treatment in turn of the eight possibilities (76), (77) proves the lemma. $`\mathrm{}`$
###### Lemma 7.19.
If a super-word $`W`$ equals one of the super-letters $`(\text{68})`$ or $`[v]^h,`$ $`[v]B,`$ $`h1,`$ then its constitution does not equal the constitution of any word in less then $`W`$ super-letters from $`B.`$
Proof. The proof is akin to Lemma 7.5 with the following tableaux:
(78)
$$\begin{array}{cc}[u_{km}]^h,[u_{km}x_{k+1}],[u_{km}u_{km+1}]& \text{deg}_k(u)\text{deg}_{m+1}(u);\\ [v_{km}]^h,[v_{km}x_{k+1}],[v_{km}v_{km1}]& 2\text{deg}_k(u)\text{deg}_{m1}(u);\\ [v_{kk+1}x_{k+2}]& \text{deg}_k(u)=0;\\ [v_k]^h& \text{deg}_k(u)\text{deg}_n(u);\\ [u_{kn1}^3x_n]& \text{deg}_k(u)2\text{deg}_n(u).\end{array}$$
$`\mathrm{}`$
###### Lemma 7.20.
If $`y=x_i,m1i>k`$ or $`y=x_i^2,m1=i>k`$ then
(79)
$$v_{km}y_{k+1}0.$$
Proof. For $`i<m1,`$ we may transpose $`y`$ by means of (35) to the left across $`x_n^2`$ and then use Lemma 37 with $`m^{}=n1.`$
If $`y=x_i^2,`$ $`m1=i>k`$ then by the above case, $`i<m1,`$ we get
(80)
$$v_{km}y=v_{km+1}\underset{¯}{x_mx_{m1}^2}=\underset{¯}{v_{km+1}x_{m1}}(\alpha x_mx_{m1}+\beta x_{m1}x_m)_{k+1}0,$$
where by definition $`v_{kn}=u_{kn}`$ and $`u_{kn}x_{n2}_{n2}0,`$ while $`n2=i>k`$.
If $`y=x_i,`$ $`i=m>k`$ then for $`m=n1`$ we may use the inequality (74), while for $`m<n1`$ we have $`v_{km}y=v_{km+1}y_1`$ where $`y_1=x_m^2.`$ Hence we may use (80) replacing $`m`$ by $`m+1.`$
If $`y=x_i,i>m>k`$ then by (35) we get $`v_{km}y=\alpha v_{ki+1}\underset{¯}{x_ix_{i1}x_i}w.`$ Changing the underlined by (33), we may apply the previously considered cases: $`m^{}1=i^{},`$ where $`m^{}=i+1,`$ $`i^{}=i;`$ and $`i^{}<m^{}1,`$ where $`m^{}=i+1,`$ $`i^{}=i1.`$ $`\mathrm{}`$
If we multiply (69) by $`x_n`$ from the right and subtract (33) with $`i=n1`$ multiplied from the left by $`x_{n1}^2,`$ then by means of $`p_{n1n1}^2=p_{nn1}p_{n1n}=p_{nn}^1`$ we get
(81)
$$x_{n1}^2\underset{¯}{x_nx_{n1}x_n}_np_{n1n}(p_{n1n1}^{[3]}x_{n1}x_nx_{n1}^2x_np_{n1n1}x_{n1}^2x_n^2x_{n1}).$$
Let us first multiply this relation by $`x_{n2}^2`$ from the left and then apply (33) to the underlined sub-word. Taking into account the relation $`x_{n2}^2x_{n1}^3_{n1}0,`$ we get that the left hand side of the multiplied (81) equals $`p_{n1n}p_{nn}(1+p_{nn})^1x_{n2}^2x_{n1}^2x_n^2x_{n1}`$ up to $`_{n1},`$ i.e. it is proportional to the second term of the right hand side. As a result the relation below with $`\alpha =p_{n1n1}^1(1+p_{nn})0`$ is correct.
(82)
$$x_{n2}^2x_{n1}^2x_n^2x_{n1}_{n1}\alpha x_{n2}^2x_{n1}x_nx_{n1}^2x_n.$$
###### Lemma 7.21.
If $`k<s<mn`$ and as above $`v_{kn}=u_{kn}`$ then
(83)
$$v_{km}v_{ks}_{k+1}\epsilon v_{ks}v_{km},\epsilon 0,$$
Proof. Let us use downward induction on $`k.`$ For this we first transpose the second letter $`x_k`$ of $`v_{km}v_{ks}`$ as far to the left as possible by means of (35), and then change the onset $`x_kx_{k+1}x_k`$ according to (72). We get
(84)
$$v_{km}v_{ks}_{k+1}\alpha x_k^2(v_{k+1m}v_{k+1s}),\alpha 0.$$
For $`k+1<s`$ we may apply the inductive supposition to the word in the parentheses, and then transpose $`x_k`$ to its former position by (72), (35).
For $`k+1=s`$ we will use downward induction on $`s.`$
Let $`k+1=s=n1.`$ In this case $`m=n`$ and (84) becomes:
$$v_{n2n}v_{n2n1}_{n1}\beta x_{n2}^2(x_{n1}\underset{¯}{x_nx_{n1}x_n}x_{n1}).$$
Let us replace the underlined part according to (33). Since $`x_{n2}^2x_{n1}x_n^2_n0,`$ we may continue by (82):
$$_{n1}\beta _1x_{n2}^2x_{n1}^2x_n^2x_{n1}_{n1}\beta _2\underset{¯}{x_{n2}^2x_{n1}}x_nx_{n1}^2x_n_{n1}$$
$$\beta _3x_{n2}x_{n1}\underset{¯}{x_{n2}x_n}x_{n1}^2x_n_{n1}\beta _4x_{n2}x_{n1}x_n\underset{¯}{x_{n2}x_{n1}^2}x_n.$$
With the help of (33) we get
$$=\epsilon v_{n2n1}v_{n2n}+\beta _5x_{n2}\underset{¯}{x_{n1}x_nx_{n1}^2}x_{n2}x_n,\epsilon 0.$$
By (73) and (71) we see that the second term equals zero up to $`_{n1}.`$
The inductive step on $`s`$ coincides the inductive step on $`s`$ in Lemma 54 up to replacing both the citations of Lemma 51 with the citations of Lemma 79 and $`w`$ with $`v.`$ $`\mathrm{}`$
###### Lemma 7.22.
The set $`B`$ satisfies the Lemma 4.8 conditions.
Proof. According to the Super-letter Crystallisation Lemma and Lemma 7.11 it is sufficient to show that words of the form $`u_0,u_1,u_2,w_3,w_4,w_5,w_6`$ are linear combinations of lesser words in $`U_P(𝔤).`$ The words $`u_0`$ are diminished by (35). The words $`u_1,u_2`$ have been diminished in Theorem $`A_n`$ since in the case $`C_n`$ the words $`u_2`$ are independent of $`x_n,`$ while $`u_1`$ depends on $`x_n`$ only if $`u_1=x_{n1}x_n^2.`$ The relation (79) shows that $`w_3_{k+1}0,`$ $`w_4_{k+1}0.`$ Lemma 83 with $`s=m1`$ gives the required representation for $`u_5.`$
Consider the words $`w_6.`$ For $`k=n1`$ the relation (69) defines the required decomposition. Let $`k<n1.`$ Since $`x_1,\mathrm{},x_{n1}`$ generate a sub-algebra of the type $`A_{n1},`$ the crystal decomposition of $`u_{kn2}^3x_{n1}`$ has the form
(85)
$$u_{kn2}^3x_{n1}=\alpha u_{m_1s_1}u_{m_2s_2}\mathrm{}u_{m_ts_t},$$
where $`u_{m_1s_1}u_{m_2s_2}\mathrm{}u_{m_ts_t},`$ that is $`m_1m_2\mathrm{}m_t`$ and $`s_is_{i+1}`$ if $`m_i=m_{i+1}.`$ In particular, if $`m_1=k`$ then $`m_2=\mathrm{}=m_t=k`$ and, due to the homogeneity, $`t=3,`$ $`s_1=n1,`$ $`s_2=s_3=n2.`$ Therefore
(86)
$$u_{kn2}^3x_{n1}_{k+1}\epsilon u_{kn1}u_{kn2}^2.$$
Along similar lines, the following relations are valid as well
(87)
$$u_{kn2}^3x_{n1}^2_{k+1}\mu u_{kn1}^2u_{kn2},u_{kn2}^2x_{n1}^3_{k+1}0.$$
Now let us multiply (33) with $`i=n2`$ by $`x_{n1}`$ from the right, and then add to the result the same relation multiplied by $`p_{n2n1}(1+p_{n1n1})x_{n1}`$ from the left. We get the following relation with $`\alpha =p_{n2n1}^2p_{n1n1}^{[3]}0.`$
(88)
$$x_{n2}x_{n1}^3=\alpha x_{n1}^2x_{n2}x_{n1}+\beta x_{n1}^3x_{n2},$$
Further, we may write
(89)
$$u_{kn1}^3=\beta _1u_{kn2}u_{kn3}\underset{¯}{x_{n1}x_{n2}x_{n1}}u_{kn1},\beta _10,$$
where for $`k=n2`$ the term $`u_{kn3}`$ is absent. Let us apply (33) with $`i=n2`$ to the underlined word. Since $`u_{kn2}u_{kn3}x_{n1}^2_{n1}0,`$ we have got
(90)
$$u_{kn1}^3_{n1}\beta _2u_{kn2}^2u_{kn3}\underset{¯}{x_{n1}^2x_{n2}x_{n1}}.$$
Let us apply (88). Taking into account the second of (87) we get
(91)
$$u_{kn1}^3_{k+1}\beta _3u_{kn2}^3x_{n1}^3.$$
Let us multiply this relation from the right by $`x_n.`$ By (69) we have
(92)
$$u_{kn1}^3x_n_{k+1}\alpha \underset{¯}{u_{kn2}^3x_{n1}}x_nx_{n1}^2+\beta \underset{¯}{u_{kn2}^3x_{n1}^2}x_nx_{n1}.$$
By means of (86) and (87) we have got
$$u_{kn1}^3x_n_{k+1}\alpha _1u_{kn1}x_nu_{kn2}^2x_{n1}^2+\beta _1u_{kn1}^2x_nu_{kn2}x_{n1},$$
and both of these words are less than $`u_{kn1}^3x_n.`$ $`\mathrm{}`$
###### Lemma 7.23.
If $`p_{11}1`$ then the values of $`[v]^h,`$ where $`[v]B,`$ $`vx_i,`$ $`h1`$ are not skew primitive. In particular they are non-zero.
Proof. Note that for $`n>3`$ the algebra generated by $`x_2,\mathrm{}x_n`$ is a sub-algebra of the type $`C_{n1}.`$ Therefore we may use induction on $`n`$ with additional supposition that the theorem statements 1 and 2 are valid for the lesser values of $`n.`$ We will formally consider the sub-algebra generated by $`x_{n1},x_n`$ as an algebra of the type $`C_2,`$ and the sub-algebra generated by $`x_n`$ as an algebra of type $`C_1.`$ In this case for $`n=1`$ the present lemma and the statements 1 and 2 are valid in obvious way.
If the first letter $`x_k`$ of $`v`$ is less than $`x_1`$ then we may use the inductive supposition directly. If $`v=u_{1m}`$ then one may literally repeat arguments of Lemma 7.6 starting at (39).
If $`v=v_{1m}`$ and $`n>3`$ then we may repeat arguments of Lemma 7.16 starting at (63) up to replacing $`w`$ with $`v.`$ For $`n=3`$ in these arguments the formula (66) assumes the form
(93)
$$\mathrm{\Sigma }=\epsilon g_1(p(v_{13},x_2)p_{21}[x_2^2x_3]+(1p_{11}^1)[x_2x_3]x_2)x_1.$$
Therefore the left component of the tensor $`\mathrm{\Sigma }`$ is a non-zero linear combination of the basis elements. For $`n=2`$ the set $`B`$ has no elements $`v_{1m}`$ at all.
Consider the last case, $`v=v_1=[u_{1n1}^2x_n].`$ Let $`S_k`$ be the sum of all tensors of $`\mathrm{\Delta }([u_{kn}])=`$ $`u^{(1)}u^{(2)}`$ with deg$`{}_{n}{}^{}(w^{(1)})=1,`$ deg$`{}_{k}{}^{}(w^{(1)})=0,`$ $`k<n.`$ Evidently $`S_n=x_n1.`$ Let us show by downward induction on $`k`$ that $`S_k=(1p_{11}^1)g(u_{kn1})x_n[u_{kn1}]`$ at $`k<n.`$ We have
(94)
$$\mathrm{\Delta }([u_{kn}])=\mathrm{\Delta }(x_k)\mathrm{\Delta }([u_{k+1n}])p(x_k,u_{k+1n})\mathrm{\Delta }([u_{k+1n}])\mathrm{\Delta }(x_k).$$
Consequently,
(95)
$$S_k=(g_kx_k)S_{k+1}p(x_k,u_{k+1n})S_{k+1}(g_kx_k).$$
This implies the required formula since by (70) at $`k<n1`$ we have
$$p(x_k,u_{k+1n})p(x_n,x_k)=p(x_k,u_{k+1n1}),$$
while at $`k=n1`$ we have $`p(x_{n1},x_n)p(x_n,x_{n1})=p_{11}^1.`$
In a similar manner, consider the sum $`S`$ of all tensors of $`\mathrm{\Delta }([u_{kn}^2x_n])=`$ $`w^{(1)}w^{(2)}`$ with deg$`{}_{n}{}^{}(w^{(1)})=1,`$ deg$`{}_{i}{}^{}(w^{(1)})=0,`$ at $`i<n.`$
(96)
$$\mathrm{\Delta }([[u_{1n1}][u_{1n}]])=\mathrm{\Delta }([u_{1n1}])\mathrm{\Delta }([u_{1n}])p(u_{1n1},u_{1n})\mathrm{\Delta }([u_{1n}])\mathrm{\Delta }([u_{1n1}]).$$
Since we now $`S_1,`$ we may calculate $`S:`$
$`S=(g(u_{1n1})[u_{1n1}])S_1p(u_{1n1},u_{1n})S_1(g(u_{1n1})[u_{1n1}])=`$
(97) $`(1p_{11}^1)g(u_{1n1}^2)x_n(1p(u_{1n1},u_{1n})p(x_n,u_{1n1}))[u_{1n1}]^2.`$
By (70), using the bicharacter property of $`p,`$ we have
$$1p(u_{1n1},u_{1n})p(x_n,u_{1n1})=1p(u_{1n1},u_{1n1})p_{n1n}p_{nn1}=$$
$$1p_{n1n1}p_{n1n1}^2=1p_{11}^10.$$
Because of this, $`S0`$ and the sum of all tensors $`w^{(1)}w^{(2)}`$ with deg$`{}_{n}{}^{}(w^{(1)})=h,`$ deg$`{}_{k}{}^{}(w^{(1)})=0,`$ $`k<n`$ of the basis decomposition of $`\mathrm{\Delta }([v_1]^h)`$ equals $`S^h0.`$ Therefore $`[v_1]^h`$ is not skew primitive. $`\mathrm{}`$
Proof of Theorem $`C_n.`$ For the first statement it will suffice to prove that all super-letters (67) are hard in $`U_P(𝔤).`$ Since none of $`u_{km},`$ $`v_{km}`$ contains a sub-word (28), Lemma 7.1 implies that $`[u_{km}],`$ $`[v_{km}]`$ are hard.
If $`[v_k]`$ is not hard then, by the homogeneous version of Definition 4.3, its value is a polynomial in lesser hard super-letters. In line with Lemmas 7.22 and 4.8, all hard super-letters belong to $`B.`$ Therefore, by Lemma 7.19, $`[v_k]=0.`$ Since deg$`{}_{n}{}^{}(v_k)=1`$ and deg$`{}_{n1}{}^{}(v_k)=2,`$ the equality $`[v_k]=0`$ is valid in the algebra $`C^{}`$ which is defined by all relations of $`U_P(𝔤),`$ but ones of degree greater than 1 in $`x_n`$ and ones of degree greater than 2 in $`x_{n1},`$ that is in the algebra defined by (33), (34) with $`i<n1,`$ and (35). These relations do not reverse the order of $`x_{n1}`$ and $`x_n`$ in monomials since none of them has both $`x_{n1}`$ and $`x_n.`$ This implies that the sum of all monomials of $`[v_k]=[u_{kn1}][u_{kn}]`$ $`p(u_{kn1},u_{kn})[u_{kn}][u_{kn1}]`$ in which $`x_n`$ is prefixed to $`x_{n1}`$ equals zero in $`C^{},`$ that is $`[u_{kn}][u_{kn1}]=0.`$ Especially, this equality is valid in $`U_P(𝔤).`$ Since, by Theorem 4.5, the super-word $`[u_{kn}][u_{kn1}]`$ is a basis element, the first statement is proved.
If $`[v]B`$ is of finite height then, by Lemma 7.19 and the homogeneous version of Definition 4.4, we have $`[v]^h=0.`$ For $`p_{11}1`$ this contradicts Lemma 7.23. In a similar manner, according to Lemma 4.9, every skew primitive homogeneous element has the form $`[v]^h.`$ This, together with Lemma 7.23, proves the fourth statement and, for $`p_{11}1,`$ the second one too. If $`p_{11}=1`$ then according to (70) we have $`p_{ii}=p_{ij}p_{ji}=1`$ at all $`i,j.`$ In particular, the skew commutator is a quantum operation. Hence all elements of $`B`$ are skew primitive. These elements span a colour Lie algebra. Now, as in Theorem $`A_n,`$ we may use the coloured PBW theorem.
The third statement will follow from Theorem 5.2 and Lemmas 5.3, 7.18 provided we note that all super-letters (68) are zero in $`U_P(𝔤).`$ We have proved already that these super-letters are non-hard. So it remains to use first the homogeneous version of Definition 4.3 and then Lemma 7.26. $`\mathrm{}`$
Theorem $`𝐃_𝐧`$. Let $`𝔤`$ be of the type $`D_n,`$ and $`p_{ii}1,`$ $`1in.`$ Denote by $`B`$ the set of the following super-letters:
(98)
$$\begin{array}{cccc}[u_{km}]& \stackrel{df}{=}& [x_kx_{k+1}\mathrm{}x_m],& 1km<n;\\ [e_{km}]& \stackrel{df}{=}& [x_kx_{k+1}\mathrm{}x_{n2}x_nx_{n1}\mathrm{}x_m],& 1k<mn,\\ [e_{n1n}]& \stackrel{df}{=}& x_n.\end{array}$$
The statements given below are valid.
$`1.`$ The values of $`(\text{98})`$ in $`U_P(𝔤)`$ form the PBW-generators set.
$`2.`$ Each of the super-letters $`(\text{98})`$ has infinite height in $`U_P(𝔤).`$
$`3.`$ The relations $`(\text{21})`$ together with the following ones form a Groebner–Shirshov system that determines the crystal basis of $`U_P(𝔤).`$
(99)
$$\begin{array}{cccc}[u_0]& \stackrel{df}{=}& [x_kx_m]=0,& 1k<m1<n,(k,m)(n2,n);\\ [u_1]& \stackrel{df}{=}& [u_{km}x_{k+1}]=0,& 1k<m<n;\\ [u_1^{}]& \stackrel{df}{=}& [x_{n2}x_n^2]=0,& \\ [u_2]& \stackrel{df}{=}& [u_{km}u_{km+1}]=0,& 1km<n1;\\ [v_3]& \stackrel{df}{=}& [e_{km}x_{k+1}]=0,& 1k<mn,n1km2;\\ [v_4]& \stackrel{df}{=}& [e_{kk+1}x_{k+2}]=0,& 1k<n2;\\ [v_4^{}]& \stackrel{df}{=}& [e_{n3n2}x_n]=0,& \\ [v_5]& \stackrel{df}{=}& [e_{km}e_{km1}]=0,& 1k<m1n1;\\ [v_6]& \stackrel{df}{=}& [u_{km}e_{kn}]=0,& 1km<n,n2m.\end{array}$$
$`4.`$ If $`p_{11}1,`$ then the generators $`x_i,`$ their powers $`x_i^t,x_i^{tl^k},`$ such that $`p_{ii}`$ is a primitive $`t`$-th root of 1, together with the constants $`1g,`$ $`gG`$ form a basis of $`𝔤_P=L(U_P(𝔤)).`$ Here $`l=`$ char$`(𝐤).`$
$`5.`$ If $`p_{11}=1,`$ then the elements of $`B`$ and, for $`l>0,`$ their $`l^k`$-th powers together with the constants $`1g,`$ $`gG`$ form a basis of $`𝔤_P.`$
In the case $`D_n`$ the algebra $`U_P^b(𝔤)`$ can be defined by the condition that the sub-algebras $`U_{n1}`$ and $`U_n`$ generated, respectively, by $`x_1,\mathrm{},x_{n1}`$ and $`x_1,\mathrm{},x_{n2},x_{n1}^{}=x_n`$ are quantum universal enveloping algebras of the type $`A_{n1},`$ and by the only additional relation
(100)
$$[x_{n1}x_n]=0.$$
The existence conditions take up the form
$`p_{ii}=p_{nn}=p_{11},p_{i+1i}p_{ii+1}=p_{n2n}p_{nn2}=p_{11}^1,\text{ if }1i<n,`$
(101) $`p_{n1n}p_{nn1}=p_{ij}p_{ji}=1,\text{ if }ij>1\&(i,j)(n,n2).`$
###### Lemma 7.24.
The brackets in $`(\text{98})`$ are set up by the recurrence formulae
(102)
$$\begin{array}{cccc}[e_{km}]& =& [x_k[e_{k+1m}]],& \text{if}1k<m1<n,kn1;\\ [e_{kk+1}]& =& [[e_{kk+2}]x_{k+1}],& \text{if}1k<n1.\end{array}$$
Proof. It is enough to use the properties 6s, 1s, and 2s. $`\mathrm{}`$
###### Lemma 7.25.
If $`[u],[v]B,`$ then one of the statements below is correct.
$`1)[[u][v]]`$ is not a standard nonassociative word;
$`2)uv`$ contains a sub-word of one of the types $`u_0,u_1,u_1^{}u_2,v_3,v_4,v_4^{},v_5,v_6;`$
$`3)[[u][v]]B.`$
Proof. The formulae (102) coincides with (49) at $`kn1`$ up to replacing $`e`$ by $`w.`$ The inequality $`e_{km}>e_{rs}`$ is set up by the same conditions, $`k<r(k=r\&m<s),`$ as the inequality $`w_{km}>w_{rs}`$ does. Likewise $`u_{km}>e_{rs}`$ is set up by the same condition, $`kr,`$ as $`u_{km}>w_{rs}`$ does. Therefore Lemmas 7.8, 7.9, 7.10 remain valid with $`e`$ in place of $`w:`$
(103)
$$\begin{array}{ccc}[[e_{km}][e_{rs}]]& \text{ is standard }& sm>k+1=r(s<m\&r=k);\\ [[u_{km}][e_{rs}]]& \text{ is standard }& k=rk=m<r;\\ [[e_{km}][u_{rs}]]& \text{ is standard }& r=k+1<mr=k+1=m=s.\end{array}$$
By looking over all of these possibilities we get the lemma statement. $`\mathrm{}`$
###### Lemma 7.26.
If a super-word $`W`$ equals one of the super-letters $`(\text{99})`$ or $`[v]^h,`$ $`[v]B,`$ $`h1`$ then its constitution does not equal the constitution of any super-word in less than $`W`$ super-letters from $`B.`$
Proof. The proof is similar to the one of Lemma 7.5 with the tableaux
$`\begin{array}{ccccc}[u_{km}]^h,& [u_{km}x_{k+1}],& [u_{km}u_{km+1}]& & \text{deg}_k(u)\text{deg}_{m+1}(u);\\ [e_{km}]^h,& [e_{km}x_{k+1}],& [e_{km}e_{km1}],& m<n& 2\text{deg}_k(u)\text{deg}_{m1}(u);\\ [e_{kn}]^h,& [e_{kn}x_{k+1}],& [e_{kn}e_{kn1}]& & \text{deg}_k(u)\text{deg}_{m1}(u);\end{array}`$
(104) $`\begin{array}{cc}[e_{kk+1}x_{k+2}]& \text{deg}_k(u)=0;\\ [e_{n3n2}x_n]& \text{deg}_{n3}(u)=0;\\ [u_{kn2}e_{kn}]& \text{deg}_k(u)\text{deg}_{n1}(u)+\text{deg}_n(u);\\ [u_{kn1}e_{kn}]& \text{deg}_k(u)\text{deg}_n(u).\end{array}`$
###### Lemma 7.27.
If $`y=x_i,`$ $`m1i>k`$ or $`y=x_i^2,`$ $`m1=i>k`$ then
(105)
$$e_{km}y_{k+1}0.$$
Proof. If $`i<m1,`$ $`mn,`$ or $`m=n,`$ $`i<n2,`$ then with the help of (35) and (100) it is possible to permute $`y`$ to the left beyond $`x_n`$ and then to use Lemma 37 for $`U_{n1}.`$
If $`m=n,`$ $`i=n2`$ then we may use Lemma 37 for $`U_n.`$
If $`y=x_i^2,`$ $`m1=i>k`$ then for $`m<n`$ by the above case we get
(106)
$$e_{km}y=e_{km+1}x_mx_{m1}^2=\underset{¯}{e_{km+1}x_{m1}}(\alpha x_mx_{m1}+\beta x_{m1}x_m)_{k+1}0.$$
For $`m=n`$ we have $`e_{kn}x_{n1}^2=`$ $`\alpha \underset{¯}{u_{kn2}x_{n1}^2}x_n_{n1}0`$ since the underlined part belongs to $`U_{n1}.`$
If $`y=x_i,`$ $`i=m>k`$ then for $`m=n`$ we may use Lemma 37 applied to $`U_n;`$ for $`m=n1`$ we may use the same lemma applied to $`U_{n1}`$ provided that beforehand we permute $`x_n`$ with $`y`$ by (100); for $`m<n1`$ we may first rewrite $`e_{km}y=e_{km+1}y_1,`$ where $`y_1=x_m^2,`$ and then use (106) with $`m+1`$ in place of $`m.`$
If $`y=x_i,`$ $`i>m>k`$ then for $`i<n`$ we have $`e_{km}y=\alpha e_{ki+1}\underset{¯}{x_ix_{i1}x_i}v.`$ Replacing the underlined word by (33) in $`U_{n1},`$ we may use the previously considered cases: $`m^{}1=i^{},`$ where $`m^{}=i+1,`$ $`i^{}=i;`$ and $`i^{}<m^{}1,`$ where $`m^{}=i+1,`$ $`i^{}=i1.`$ For $`i=n,`$ and $`m=n1`$ we have $`e_{kn1}x_n=\alpha \underset{¯}{u_{kn2}x_n^2}x_{n1}`$ and one may apply Lemma 37 to $`U_n.`$ Finally, for $`i=n`$ and $`m<n1`$ we get
$$e_{km}x_n=\beta _1u_{kn2}x_nx_{n1}x_{n2}x_nv=\beta _2u_{kn2}x_{n1}\underset{¯}{x_nx_{n1}x_n}v=$$
$$\beta _3\underset{¯}{u_{kn2}x_{n1}x_{n2}}x_n^2v+\beta _4u_{kn2}\underset{¯}{x_{n1}x_n^2}x_{n2}v.$$
One may apply first Lemma 37 for $`U_{n1}`$ to the underlined sub-word of the first term, and then, after (100), Lemma 37 for $`U_n`$ to the second term. $`\mathrm{}`$
###### Lemma 7.28.
If $`k<s<mn`$ then $`e_{km}e_{ks}_{k+1}\epsilon e_{ks}e_{km},`$ $`\epsilon 0.`$
Proof. Let us carry out downward induction on $`k.`$ The largest value of $`k`$ equals $`n2.`$ In this case $`s=n1,`$ $`m=n`$ and we have
$`\underset{¯}{x_{n2}x_nx_{n2}}x_nx_{n1}_nx_{n2}^2\underset{¯}{x_n^2x_{n1}}=\alpha \underset{¯}{x_{n2}^2x_{n1}}x_n^2_{n1}`$
(107) $`\beta x_{n2}x_{n1}\underset{¯}{x_{n2}x_n^2}_n\epsilon x_{n1}x_nx_{n1}x_{n2}x_n.`$
Let us first transpose the second letter $`x_k`$ of $`e_{km}e_{ks}`$ as far to the left as possible by (35), and then replace the onset $`x_kx_{k+1}x_k`$ by (36). We get
(108)
$$e_{km}e_{ks}_{k+1}\alpha x_k^2(e_{k+1m}e_{k+1s}),\alpha 0.$$
For $`k+1<s`$ it suffices to apply the inductive supposition to the word in the parentheses and then by (36) and (35) to put $`x_k`$ to the proper place.
For $`k+1=s`$ one may use downward induction on $`s.`$ The basis of this induction, $`s=n1,`$ has been proved, see (107). For $`k<n3`$ the inductive step on $`s`$ coincides with the one of Lemma 54 with $`e`$ in place of $`w`$ since in this case the active variables $`x_k,`$ $`x_{k+1}`$ $`q`$-commute with $`x_n.`$ If $`k=n3`$ then in consideration of Lemma 54 the variable $`x_{k+1}=x_{n2}`$ is transposed across $`x_n`$ twice: in (58) and in the second word of (60).
In (58) with $`k=n3`$ we have $`s=n2,`$ $`m=n;`$ and (58) becomes
(109)
$$e_{n3n}e_{n3n2}_{n2}\beta e_{n3n1}\underset{¯}{x_{n3}x_{n2}x_nx_{n2}}.$$
In view of Lemma 105, we may transform the underlined part in $`U_n`$ neglecting the words starting with $`x_{n2}^2`$ and $`x_n`$ in much the same way as in (59), with $`x_n`$ in place of $`x_{k+1}`$. So (109) reduces to the required form.
The second word of (60) with $`k=n3`$ assumes the form $`e_{n3n}^2x_{n2}x_{n1}^2=e_{n3n}\underset{¯}{x_{n3}x_{n2}x_nx_{n2}}x_{n1}^2.`$ By Lemma 37 applied to $`U_n,`$ the underlined word is a linear combination of words starting with $`x_{n2}`$ and $`x_n.`$ However, by Lemma 105 both $`e_{n3n}x_{n2}`$ and $`e_{n3n}x_n`$ equal zero up to $`_{n2}.`$ $`\mathrm{}`$
###### Lemma 7.29.
The set $`B`$ satisfies the conditions of Lemma $`\text{4.8}.`$
Proof. By Lemmas 7.25 and 4.7 one need show only that in $`U_P^b(𝔤)`$ the words (99) are linear combinations of lesser ones. The words $`v_6`$ with $`m=n2,`$ and $`u_0,`$ $`u_1,`$ $`u_1^{},`$ $`u_2`$ have the required decomposition since they belong either to $`U_{n1}`$ or to $`U_n.`$ Lemma 105 shows that $`v_3_{k+1}0,`$ $`v_4_{k+1}0,`$ $`v_4^{}_{k+1}0.`$ Lemma 7.28 with $`s=m1`$ yields the required representation for $`v_5.`$ Consider $`v_6`$ with $`m=n1.`$ Let us prove by downward induction on $`k`$ that
$$u_{kn1}e_{kn}_{k+1}\epsilon e_{kn}u_{kn1},\epsilon 0.$$
For $`k=n1`$ this equality assumes the form (100). Let $`k<n1.`$ Let us transpose the second letter $`x_k`$ of $`u_{kn1}e_{kn}`$ as far to the left as possible in $`U_{n1}.`$ After an application of (33) we get
$$u_{kn1}e_{kn}_{k+1}\alpha x_k^2(u_{k+1n1}e_{k+1n}),\alpha 0.$$
It suffices to apply the inductive supposition to the term in the parentheses, and then by (33) and (35) for $`U_n`$ to move $`x_k`$ to the proper place. $`\mathrm{}`$
###### Lemma 7.30.
If $`p_{11}1`$ then the values of $`[v]^h,`$ where $`[v]B,`$ $`vx_i,`$ $`h1`$ are not skew primitive, in particular they are non-zero.
Proof. One need consider only super-letters that belong neither to $`U_{n1}`$ nor to $`U_n.`$ That is $`[e_{km}]`$ with $`m<n.`$ We use induction on $`n.`$
For $`n=3`$ the algebra of the type $`D_3`$ reduces to the algebra of the type $`A_3`$ with a new ordering of variables $`x_2>x_1>x_3.`$ Therefore we may use Theorem $`A_n,`$ after the decomposition below of $`e_{12}`$ in the PBW-basis:
$$[[x_1x_3]x_2]=p_{12}p_{32}[x_2[x_1x_3]]+\beta [x_1x_3]x_2.$$
Let $`n>3.`$ If $`k>1`$ then the inductive supposition works. For $`k=1,`$ $`m>2`$ we have $`e_{1m}=[x_1[e_{2m}]],`$ and one may repeat the arguments of Lemma 7.6 with $`e`$ in place of $`u`$ starting at (39). If $`m=2`$ then we may repeat the arguments of Lemma 7.16 with $`e`$ on place of $`w`$ starting at (63). $`\mathrm{}`$
Proof of Theorem $`D_n.`$ For the first statement it will suffice to prove that all super-letters (98) are hard in $`U_P^b(𝔤).`$
Since none of $`u_{km}`$ contains sub-words (28), $`[u_{km}]`$ are hard.
Suppose $`[e_{km}]`$ is non-hard. By Lemmas 7.29 and 4.8 all hard super-letters belong to $`B.`$ Thus, by Lemma 7.26, we get $`[e_{km}]=0.`$ Since deg$`{}_{n}{}^{}(e_{km})=`$deg$`{}_{n1}{}^{}(e_{km})=1,`$ the equality $`[e_{km}]=0`$ is also valid in the algebra $`D^{}`$ defined by the same relations as $`U_P^b(𝔤)`$ is, but $`[x_{n2}x_n^2]=0`$ and $`[x_{n2}x_{n1}^2]=0.`$ Let us equate to zero all monomials in all the defining relations of $`D^{},`$ but $`[x_{n1}x_n]=0.`$ Consider the algebra $`R^{}`$ defined by (100) and by the resulting system of monomial relations. It is easy to verify that the mentioned relations system $`\mathrm{\Sigma }`$ of $`R^{}`$ is closed under the compositions. Since $`e_{km}`$ contains none of leading words of $`\mathrm{\Sigma },`$ the super-letter $`[e_{km}]`$ is non-zero in $`R^{},`$ and so in $`D^{}`$ too. This contradiction proves the first statement.
If $`[v]^h,[v]B`$ is of finite height then by Lemma 7.26 and the homogeneous version of Definition 4.4 we have $`[v]^h=0.`$ For $`p_{11}1`$ this contradicts Lemma 7.30. In a similar manner, by Lemma 4.9, every skew primitive homogeneous element has the form $`[v]^h.`$ This, together with Lemma 7.30, proves both the fourth statement and the second one with $`p_{11}1.`$
If $`p_{11}=1`$ then by (101) we have $`p_{ii}=p_{ij}p_{ji}=1`$ for all $`i,j.`$ This means that the skew commutator itself is a quantum operation. Hence all elements of $`B`$ are skew-primitive. These elements span a colour Lie super-algebra. Now, as in Theorem $`A_n,`$ one may use the PBW theorem for colour Lie super-algebras.
For the third statement it will suffice to show that all super-letters (99) are zero in $`U_P(𝔤).`$ We have proved already that they are non-hard. Therefore it remains to use the homogeneous version of Definition 4.3 and Lemma 7.26. $`\mathrm{}`$
8. Conclusion
We see that in all Theorems $`A_n`$$`D_n`$ the lists of hard super-letters are independent of the parameters $`p_{ij}.`$ This fact signifies that the Lalonde–Ram basis of the ground Lie algebra (see, \[26, Figure 1\]) with the skew commutator in place of the Lie operation coincides with the set of all hard super-letters. It is very interesting to clarify how general this statement is. On the one hand, this does not hold without exception for all quantum enveloping algebras since in Theorems $`A_n`$$`D_n`$ a restriction does exist. If $`p_{ii}=1,`$ $`1i<n,`$ $`n>2`$ then it is easy to see by means of the Diamond Lemma that the sets of hard super-letters are infinite. On the other hand, this is not a specific property of Lie algebras defined by the Serre relations. By the Shirshov theorem any relation can be reduced to a linear combination of standard nonassociative words.
###### Corollary 8.1.
If $`𝔤`$ is defined by the only relation $`f=0,`$ where $`f`$ is a linear combination of standard nonassociative words, then the set of all hard in $`U_P(𝔤)`$ super-letters coincides with the Hall–Shirshov basis of $`𝔤`$ with the skew commutator in place of the Lie operation.
Proof. The only relation $`f^{}=0`$ forms a Groebner–Shirshov system since, according to 1s, none of onsets of its leading word, say $`w,`$ coincides with a proper terminal of $`w`$. Consequently, a super-letter $`[u]`$ is hard if and only if $`u`$ does not contain $`w`$ as a sub-word. We see that this criteria is independent of $`p_{ij}`$ as well. $`\mathrm{}`$
Furthermore, the third statement of Theorem $`A_n`$ shows that $`U_P^b(𝔤)`$ can be defined by the following relations in the PBW-generators $`X_u=[u].`$
(110)
$$\begin{array}{ccc}[X_u,X_v]=0,& u>v,& [[u][v]]B\\ [X_u,X_v]=X_{uv},& & [[u][v]]B.\end{array}$$
This is an argument in favour of considering the super-letters PBW-generators k$`[G]`$-module as a quantum analogue of a Lie algebra. However in the cases $`B_n,`$ $`C_n,`$ $`D_n`$ the defining relations became more complicated. For example,
(111)
$$\begin{array}{cc}B_n:[[u_{kn1}][w_{kn}]]=\alpha [u_{kn}]^2,& \alpha 0\text{ if }p_{nn}1;\\ C_n:[[u_{kn2}][v_{kn1}]]=\alpha [v_k]+\beta [u_{kn}][u_{kn1}],& \beta 0\text{ if }p_{11}1;\\ D_n:[[u_{kn2}][e_{kn1}]]=\alpha [e_{kn}][u_{kn1}],& \alpha 0\text{ if }p_{11}\pm 1.\end{array}$$
It is far more interesting that for $`p_{11}1`$ the algebra $`𝔤_P`$ turns out to be very simple in structure. Only unary quantum operations can be non-zero. Other ones may be defined, but due to the homogeneity their values equal zero. In particular, if $`p_{11}^{[t]}0`$ then without exception all quantum operations have zero values. This provides reason enough to consider $`U_P(𝔤)=`$ $`U(𝔤_P)`$ as an algebra of ‘commutative’ quantum polynomials. Certainly it is very interesting to elucidate to what extent this statement is still retained for the quantum universal enveloping algebras of homogeneous components of other Kac–Moody algebras defined by the Gabber–Kac relations (11). Also it is interesting to investigate the structure of other ‘commutative’ quantum polynomial algebras. For example, one may note that if a semi-group generated by $`p_{ij}p_{ji}`$ does not contain 1, then $`Gx_1,\mathrm{},x_n`$ itself is a ‘commutative’ quantum polynomial algebra merely since in this case there exists no non-zero quantum operation at all. In another extreme case when $`p_{ij}p_{ji}=1`$ for all $`i,j,`$ the ‘commutative’ quantum variables commute by $`x_ix_j=p_{ij}x_jx_i.`$
In a similar manner, the Drinfeld–Jimbo enveloping algebra can be considered as a ‘quantum’ Weyl algebra of (skew) differential operators (see Sec. 6). The resulting ‘quantum’ Weyl algebra is simple in the following sense.
###### Corollary 8.2.
Let $`𝔤`$ be a simple finite dimensional Lie algebra of the infinite series. If $`q^{[m]}0,`$ $`m2`$ then every non-zero Hopf ideal $`I`$ of the Drinfeld–Jimbo enveloping algebra contains all generators $`x_i,x_i^{}.`$
Proof. By the Heyneman–Radford theorem, the ideal $`I`$ has a non-zero skew primitive element, say $`a.`$ According to Lemma 6.2 and Theorems $`A_n`$$`C_n,`$ the element $`a`$ is either a constant, $`\alpha (1g),`$ or proportional to one of the elements $`x_i,x_i^{}.`$ In the former case $`I`$ contains all $`x_i`$ with $`\chi ^i(g)1`$ since $`x_ia\chi ^i(g)a_ix_i=\alpha (1\chi ^i(g))x_i.`$ Here the equality $`\chi ^i(g)=1`$ can not be valid for all $`i`$ since $`\chi ^i(g_j)=q^{d_ia_{ij}}`$ (see, Example 4 of Section 2) and the columns of the Cartan matrix are linearly independent. In the latter case (and now in the former one as well) we get $`[x_i,x_i^{}]=\epsilon _i(1g_i^2)I,`$ i.e. as above $`I`$ contains all elements $`y=x_i^\pm `$ with $`1\chi ^y(g_i^2)`$ $`=q^{\pm 2d_ja_{ij}}.`$ Since the Coxeter graph is connected, $`I`$ contains all $`x_i,x_i^{}.`$ $`\mathrm{}`$
Acknowledgements. The author is grateful to Dr. J. A. Montaraz, the director of the FES-C UNAM, Dra. S. Rodríguez-Romo, and A.V. Lara Sagahon for providing facilities for the research and also to Dr. L.A. Bokut’ and Dr. R. Bautista for helpful comments on the subject matter.
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warning/0002/quant-ph0002092.html | ar5iv | text | # Fast quantum gates for cold trapped ions.
## I Introduction
The last few years have seen impressive progress in the experimental demonstration of quantum information processing . Among the growing number of possible physical scenarios for these demonstrations, the system of laser-cooled trapped ions still remains one of the most experimentally attractive (for reviews of ion-trap quantum computing, see e.g. ). Ever since the original ion-trap proposal of Cirac and Zoller (CZ) , a number of modifications and extensions to their idea have been proposed . Many of these have aimed at bypassing two experimental hurdles of CZ’s proposal, namely: a) cooling the ionic motion to the ground state, while b) at the same time keeping the ions sufficiently far apart that individual laser access to each of them is possible. On the one hand, ‘hot’ gate implementations have been suggested that aim to function even in the presence of moderate motional heating. On the other hand, an ingenious method has been suggested that exploits the ionic micromotion induced by dc offset potentials to address individual ions even while simultaneously illuminating all ions with the same beam . Each of these proposals has its own merits and difficulties, and their feasibility and/or scalability have yet to be demonstrated in experiment. In the meantime, at least one experiment currently under development aims to tackle the two problems directly, achieving the conditions required by CZ.
In the present paper we assume that these conditions will indeed become feasible, and focus instead on another aspect of these experiments: the gate switching rates. Evidently, it is desirable that these should be as large as possible, so that a reasonably complex sequence of quantum operations can be realised before decoherence sets in. It has been remarked , that the speed of any 2-qubit gate realised by coupling two ions via a motional mode must be bounded from above by the frequency of that mode (roughly speaking, the ions must be able to ‘realise that they are moving’ before they can influence each other). At the moment of writing, no experiment realising a true 2-qubit ion-ion gate has been reported. However, at least two experiments have used schemes similar to the CZ proposal to implement 2-qubit gates between a single ion and a motional mode . Strikingly, in both cases the reported gate speeds fell far short of the mode frequency, by two to three orders of magnitude. This limitation was not circumstantial, but inherent in the experimental technique that was used. The problem was the existence of strong off-resonant ion-mode transitions, whose unwanted driving would spoil the desired gate dynamics . In order to avoid this, the laser power had to be kept at a relatively modest level, resulting in slow gates. (Very recently, a modification of the CZ scheme which allows somewhat faster gates has been proposed , see endnote).
In this paper, we propose a new scheme for 2-qubit gates which should allow an increase in gate speed by at least an order of magnitude with respect to these experiments. Furthermore, this gain is achieved without significant changes in experimental requirements with respect to existing setups, apart from an increase in laser power and good intensity stability. The key feature of our scheme is that it exploits the AC Stark-shift (lightshift) induced by light resonant with the ionic carrier transition. Using a coordinate transformation suggested by Moya-Cessa and co-workers we demonstrate that, within the Lamb-Dicke regime, and at specific shift magnitudes (i.e, laser intensities), the ion-mode dynamics assumes the form of a Jaynes-Cummings interaction . This interaction can be exploited to generate a 2-qubit gate in a manner analogous to the CZ proposal.
We then proceed to compare our scheme with other existing proposals for faster cold-atom gates. For example, already in it has been pointed out that if the travelling-wave radiation used in current experiments is replaced with a standing laser field, with the ion located at a node, then a substantial increase in gate speed would be possible. The elegant ‘Magic Lamb-Dicke parameter’ (MLDP) method proposed by Monroe et al could also in principle lead to faster gates. We argue however that our method, or possibly a combination of it with the MLDP method, is the one most amenable to practical implementation within the cold-ion scenario.
The paper is organized as follows: in the first section, we introduce our gate scheme, explaining its basic principle, the pulse sequences it requires and the ways in which it differs from existing schemes. We also discuss its scalability to many-atom arrays. We then provide numerical confirmation of our analysis, and compare the performance of our scheme with that of the Cirac-Zoller scheme in both its regimes (using travelling- or standing-wave radiation). Finally, we present our conclusions.
## II 2-qubit gates based on the AC Stark-shift effect
An important feature of the Cirac-Zoller gate scheme is that the frequencies of the pulses it uses are chosen to be resonant with the transitions between the ‘bare’ (uncoupled) ion-mode levels. This choice reflects a ‘perturbative’ point of view in which these level spacings are assumed to be unaffected by the coupling itself, or in other words that the level shifts due to the AC potential of the coupling field itself can be disregarded. For a sufficiently strong field, this assumption breaks down and the normally disregarded off-resonant transitions become important (see e.g. , sec. 4.4.6). A number of authors have speculated that it might be possible to design a gate scheme incorporating these shifts as an integral feature . In this section we construct a concrete realisation of this idea, implementing 2-bit gates by exploiting the lightshift generated by light resonant with the ionic carrier.
### A One ion interacting with a travelling laser field
In order to present our underlying idea in its clearest form, we consider first the relatively simple situation of a single trapped ion interacting with a travelling-wave field. Also for simplicity, we assume the relevant ionic levels to be coupled by a direct (optical) transition. As is well-known, the analysis can be straightforwardly adapted to the case of a Raman two-photon transition by a suitable redefinition of parameters . In later sections we demonstrate how the scheme is scalable to traps containing an $`N`$-ion chain, allowing 2-qubit gates to be realised between the internal states of any two of the ions.
In the standard interaction representation, the Hamiltonian for the one-ion system can be written as
$$H=\mathrm{}\mathrm{\Omega }[\sigma _+\mathrm{exp}(i\eta [ae^{i\nu t}+a^{}e^{i\nu t}]i\delta t)+h.c.].$$
(1)
Here, $`\delta =\omega _l\omega _a`$ is the laser-atom detuning, $`\nu `$ the trap frequency, $`\eta =\sqrt{\frac{\mathrm{}k^2}{2m\nu }}`$ is the Lamb-Dicke parameter of the trap and we have already taken into account a rotating-wave approximation (RWA) that assumes $`\delta \omega _a+\omega _l`$ (the detuning is far smaller than optical frequencies).
Let us briefly recapitulate the approach that is usually taken to this problem (see, e.g. and references therein for detailed treatments). First, one expands the exponentials in powers of $`a,a^{}`$ and looks for the resonances that arise whenever the laser frequency is tuned to a motional sideband, i.e., $`\delta =\pm m\nu `$. A second RWA is then realised, ignoring off-resonant terms which rotate at multiples of the trap frequency $`\nu `$. The remaining resonant terms can be interpreted in general as intensity-dependent ‘multiphonon’ transitions . If the Lamb-Dicke parameter is also small $`\left(\eta 1\right)`$, and the ion is sufficiently cooled, the intensity dependence of the coupling constant can be ignored to lowest order in $`\eta `$. For example, if the laser is resonant with the carrier transition ($`m=0`$), or with the first red sideband ($`m=1`$), we have respectively the simple forms
$`H_{1CZ}`$ $``$ $`\mathrm{}\mathrm{\Omega }e^{\frac{1}{2}\eta ^2}\left[\sigma _++\sigma _{}\right]`$ (3)
$`H_{2CZ}`$ $``$ $`i\mathrm{}\mathrm{\Omega }\eta e^{\frac{1}{2}\eta ^2}\left[\sigma _+a\sigma _{}a^{}\right].`$ (4)
These are the interactions that form the basis of the standard Cirac-Zoller scheme for realising 1- and 2-qubit quantum logic gates . A slight modification of this scheme (using blue-sideband-detuned pulses) has been implemented experimentally in single-ion traps .
### B 2-qubit lightshift-based quantum gates
We now demonstrate that, even if only radiation resonant with the carrier is used, and without leaving the Lamb-Dicke limit, there is still a regime where 2-qubit dynamics can be obtained. The basic physical idea behind this is as follows: we know that, apart from driving the 1-qubit transition described in Eq.(3), any radiation resonant with the carrier will also lead to an AC level-splitting of the ionic semiclassical dressed states $`|\pm =\frac{1}{\sqrt{2}}\left(|g\pm |e\right)`$ (Fig. 1$`a`$) ). The magnitude of the splitting is $`2\mathrm{}\mathrm{\Omega }`$, where $`\mathrm{\Omega }`$, is the Rabi frequency. When the intensity of the laser is such that the splitting equals exactly one vibrational energy quantum $`\mathrm{}v`$, the levels $`|+|0`$ and $`||1`$ become degenerate, and we can expect transitions between them (Fig. 1$`b`$). This amounts effectively to an exchange of excitation between the motional and internal states, i.e., to 2-qubit dynamics.
To see how this happens in detail, let us begin by first making the Lamb-Dicke approximation (to first order in $`\eta `$) directly in eq. $`\left(\text{1}\right)`$
$`H`$ $``$ $`\mathrm{}\mathrm{\Omega }e^{\frac{1}{2}\eta ^2}[\sigma _+e^{i\delta t}(1+i\eta [ae^{i\nu t}+a^{}e^{i\nu t}])+h.c.]`$ (5)
$`=`$ $`\mathrm{}\mathrm{\Omega }^{}\left[\begin{array}{c}\left(\sigma _+e^{i\delta t}+\sigma _{}e^{+i\delta t}\right)+\\ i\eta \left(\sigma _+e^{i\delta t}\sigma _{}e^{+i\delta t}\right)\left[ae^{i\nu t}+a^{}e^{i\nu t}\right]\end{array}\right]`$ (8)
(where we have defined $`\mathrm{\Omega }^{}\mathrm{\Omega }e^{\frac{1}{2}\eta ^2}`$). When the radiation is resonant with the ionic transition (or ‘carrier’) frequency $`\left(\delta =0\right)`$, this reduces to
$$H\mathrm{}\mathrm{\Omega }^{}\left[\sigma _++\sigma _{}+i\eta \left(\sigma _+\sigma _{}\right)\left[ae^{i\nu t}+a^{}e^{i\nu t}\right]\right].$$
(9)
Comparing with eq. (3), we see that the usual derivation corresponds to neglecting the terms rotating at frequency $`\pm \nu `$ in this expression. These terms are the first order correction to the semiclassical ion-field interaction due to the presence of the trapping potential, and their effect is to cause the dressed states $`|\pm `$ to become nonstationary. To see how these evolve, we first move into the ‘dressed-state’ picture obtained by rotating the atomic basis states with the transformation
$$R=\frac{1}{\sqrt{2}}\left(\begin{array}{cc}1\hfill & 1\hfill \\ 1\hfill & 1\hfill \end{array}\right),$$
(10)
so that $`|\pm `$ become respectively $`|e`$ and $`|g`$ (note that, in our notation, $`|e=\left(\genfrac{}{}{0pt}{}{1}{0}\right)`$, $`|g=\left(\genfrac{}{}{0pt}{}{0}{1}\right)`$, $`\sigma _+=\left(\genfrac{}{}{0pt}{}{\mathrm{0\hspace{0.17em}1}}{\mathrm{0\hspace{0.17em}0}}\right)`$, $`\sigma _{}=\left(\genfrac{}{}{0pt}{}{\mathrm{0\hspace{0.17em}0}}{\mathrm{1\hspace{0.17em}0}}\right)`$).
Using the fact that
$$R\sigma _\pm R^{}=\frac{1}{2}\left(\sigma _z\pm \left(\sigma _+\sigma _{}\right)\right)$$
(11)
we can see that, in this picture, the Hamiltonian has the Jaynes-Cummings form
$$H^{}=\mathrm{}\mathrm{\Omega }^{}\left[\sigma _z+i\eta \left(\sigma _+\sigma _{}\right)\left[ae^{i\nu t}+a^{}e^{i\nu t}\right]\right].$$
(12)
(This transformation of the Hamiltonian is a special case of the construction given in , where it is shown that the ion-laser interaction is always unitarily equivalent to a Jaynes-Cummings form, without any approximations).
Making a further ‘interaction picture’ transformation of the Hamiltonian by the unitary operator $`\mathrm{exp}\left(\frac{i\mathrm{\Omega }^{}t\sigma _z}{\mathrm{}}\right)`$, we have
$$H^{\prime \prime }=i\mathrm{}\eta \mathrm{\Omega }^{}\left[\begin{array}{c}e^{i\left(2\mathrm{\Omega }^{}\nu \right)t}\sigma _+ae^{i\left(2\mathrm{\Omega }^{}\nu \right)t}\sigma _{}a^{}\\ +e^{i\left(2\mathrm{\Omega }^{}+\nu \right)t}\sigma _+a^{}e^{i\left(2\mathrm{\Omega }^{}+\nu \right)t}\sigma _{}a\end{array}\right],$$
(13)
which gives us the resonance condition
$$\mathrm{\Delta }=\mathrm{\Omega }^{}\frac{\nu }{2}=0.$$
(14)
Apart from the small correction to $`\mathrm{\Omega }`$ given by the Debye-Waller factor $`e^{\frac{1}{2}\eta ^2}`$ , this is precisely the condition depicted in Fig. 1. In this case, the first two (‘rotating’) terms in Eq. (13) become constant while the second two (‘counter rotating’) oscillate at frequency $`2\nu `$. We can ignore them, making the Jaynes-Cummings RWA, as long as the secular frequency $`\eta \mathrm{\Omega }^{}=\frac{1}{2}\eta \nu `$ of the resulting evolution is much smaller than this . This requires $`\eta 4`$, which is compatible with the Lamb-Dicke assumption $`\eta 1`$ we have already made. Thus, if the laser’s frequency and intensity are such that they satisfy the double resonance condition $`\delta =\mathrm{\Delta }=0`$, the evolution of the system can be described by the simple Jaynes-Cummings form
$$H_{2SS}=\frac{i\mathrm{}\eta \nu }{2}\left[\sigma _+a\sigma _{}a^{}\right].$$
(15)
What this teaches us is that off-resonant transitions cannot always be disregarded, but, under the right conditions, may in fact lead to resonant effects. Intuitively, if the off-resonant terms in the Hamiltonian given in eq. (9) rotate precisely in step with the secular evolution generated by the resonant terms, their contribution does not ‘average out’ but rather adds up over each cycle, in a manner reminiscent of an oscillator being driven by a resonant force. In the present case this effect allows a field resonant with the carrier to couple the internal and motional ionic variables in a way exactly analogous to a red-sideband-detuned pulse as described by eq. (4) (Fig. 2$`a,b`$). In particular, it can just as well be used to implement $`2`$-qubit logic gates between these two degrees of freedom. Of course, the Hamiltonian (15) is valid only in the ‘dressed’ picture defined by the operator $`R`$ in eq. (10). In the normal or ‘bare’ picture, its effect can be seen as a beating at frequency $`\eta \nu `$ superposed on the usual Rabi flops between states $`|g|n`$ and $`|e|n`$ (Fig. 2$`c`$). It is not necessarily obvious that this ‘dressed-picture’ Jaynes-Cummings interaction can be used to implement quantum logic gates in the ‘real world’. Nevertheless, in the next section we show how, with a suitable generalisation to the $`N`$-ion situation, this interaction can indeed realise a Control-NOT (C-NOT) gate between the internal variables of two separate ions. (Recall that a C-NOT gate together with one-qubit rotations form a universal set of gates for quantum computing ).
Finally, let us briefly consider the experimental requirements of our proposal. Apart from the usual demands of the CZ quantum gate proposal (individual ion access, ground-state cooling), the only new requirement we make is that the laser should have a fixed intensity satisfying the resonance condition in eq. (14) (or its $`N`$-ion generalisation, see below). In more quantitative terms, our numerical simulation (see section III B 1) indicates that the laser power must be stable to within about $`\pm 0.5\%`$. This does not seem to require significant improvements in the laser power and intensity stability already available in current experimental setups . There is also a bonus in the fact that a single laser can be used to perform both 1- and 2- qubit interactions. Therefore, we expect that a proof-of-principle experiment using a single trapped ion should not be hard to realise.
### C Lightshift gates in a chain of N ions
The results we have just described are almost immediately generalisable to the case where there are $`N`$ identical ions (and therefore $`N`$ motional modes) in a linear trap . Assuming that each of the ions can be illuminated individually by a (travelling) laser beam, then resonance conditions similar to eq. (14) turn out to exist for each separate mode frequency $`\nu _j`$. Before showing how the resulting ion-mode interaction can be used to implement ion-ion gates, we would like to call attention to an important aspect of the $`N`$-ion situation. In principle, any of the $`N`$ motional modes can be used to couple the internal ionic variables. However, in order for the lightshift scheme to function with higher-order modes it is necessary to drive the system deeper into the Lamb-Dicke regime. To see this, consider the interaction-picture Hamiltonian describing the coupling of the $`j^{th}`$ ion with a (travelling-wave) laser
$$H=\mathrm{}\mathrm{\Omega }[\sigma _+^j\mathrm{exp}(i\underset{p=1}{\overset{N}{}}\eta _{jp}[a_pe^{i\nu _pt}+a_p^{}e^{i\nu _pt}]\delta t)+h.c.]$$
(16)
Here, $`p`$ indexes the normal modes. The parameter $`\eta _{jp},`$which functions as the ‘effective’ Lamb-Dicke parameter of the $`p^{th}`$ mode, corresponds to the product $`\eta _pb_j^{\left(p\right)}`$, where $`\eta _p=\sqrt{\frac{\mathrm{}k^2}{2m\nu _p}}`$ is the ‘conventional’ Lamb-Dicke parameter and $`b_j^{\left(p\right)}`$is the relative weight of the $`j^{th}`$ ion’s displacement in this mode. For the centre-of-mass mode, $`b_j^{\left(1\right)}=\frac{1}{\sqrt{N}}`$ is independent of which ion is being driven. For other modes this is no longer true. James has given values of $`b_j^{\left(p\right)}`$ for all ions and modes up to $`N=10`$.
If all modes are suitably cooled and within the Lamb-Dicke regime, and if the laser is resonant with the ionic carrier transition $`\left(\delta =0\right)`$, then a procedure entirely analogous to the one described in eqs. (8-13) can be followed. One then obtains that, in the ‘dressed-state’ picture defined by
$$V\left(t\right)=\frac{1}{\sqrt{2}}\mathrm{exp}\left(i\mathrm{\Omega }^{}t\sigma _z^j\right)R_j,$$
(17)
the Hamiltonian given above can be rewritten as
$$H^{\prime \prime }=i\mathrm{}\mathrm{\Omega }^{}\underset{p}{}\eta _{jp}\left[\begin{array}{c}e^{i\left(2\mathrm{\Omega }^{}\nu _p\right)t}\sigma _+^ja_pe^{i\left(2\mathrm{\Omega }^{}\nu _p\right)t}\sigma _{}^ja_p^{}+\\ +e^{i\left(2\mathrm{\Omega }^{}+\nu _p\right)t}\sigma _+^ja_p^{}e^{i\left(2\mathrm{\Omega }^{}+\nu _p\right)t}\sigma _{}^ja_p\end{array}\right],$$
(18)
where $`\mathrm{\Omega }^{}=\mathrm{\Omega }e^{\frac{1}{2}\left(_p\eta _{jp}^2\right)}`$. As expected, there are multiple resonance conditions analogous to eq. (14), one for each mode frequency $`\nu _p.`$ If any of these are met (say, $`\mathrm{\Omega }^{}=\frac{\nu _q}{2}`$ for the $`q^{th}`$ mode), then the terms in this Hamiltonian can be divided into three categories according to their time-dependence: (1) The rotating terms of the $`q^{th}`$ mode are resonant, and represent a Jaynes-Cummings interaction of the form
$$H_{2SS}=\frac{i\mathrm{}\nu _q\eta _{jq}}{2}\left(\sigma _+^ja_p\sigma _{}^ja_p^{}\right);$$
(19)
(2) All counter-rotating terms oscillate at frequencies equal (in modulus) to at least $`\nu _q+\nu _1\frac{\nu _q\eta _{jq}}{2}`$, where $`\nu _1`$ is the lowest energy mode. Assuming the effective Lamb-Dicke parameter $`\eta _{jq}`$ is small $`(\eta _{jq}\frac{1}{10})`$, they can therefore can be discarded in a RWA. (3) The rotating terms of the other modes oscillate at frequencies equal to $`\pm \left|\nu _p\nu _q\right|`$. For a similar $`\eta _{jq}`$ these terms can be discarded as long as
$$\frac{\left|\nu _p\nu _q\right|}{\nu _q}\frac{\eta _{jq}}{2}.$$
(20)
If this is true for all $`pq`$, then the Hamiltonian (18) can be reduced to the resonant term given in eq. (19). In this case, only the $`q^{th}`$ mode is coupled to the ion’s internal state, just as in the usual perturbative scheme when the laser is tuned to the first red sideband of this mode. The off-resonant terms will lead to a small population leakage into the unwanted modes, of order
$$ϵ^2=\left(\frac{\eta _{jq}\nu _q}{2\left|\nu _p\nu _q\right|}\right)^21.$$
(21)
As we discuss in appendix A, for high enough precision (small enough $`ϵ`$), this population loss gives an upper bound to $`\eta _{jq}`$, and therefore to the overall Rabi frequency $`\frac{1}{2}\eta _jq\nu _q`$ at which the scheme can function. For example, in the case of the lowest (center-of-mass) mode, $`ϵ^20.005`$ requires $`\eta _10.1`$. In addition, it has been shown by James that the spacing $`|\nu _{q+1}\nu _q|`$ between successive modes decreases as their order increases. It follows that attaining a given precision $`ϵ`$ requires $`\eta _{jq}`$ to be made smaller and smaller as $`q`$ grows. In effect, we find that the potential increase in Rabi frequency afforded by using higher modes is completely counterbalanced by this requirement, with the result that the maximum value for the overall switching rate actually decreases as higher modes are used.
### D 2-ion CNOT gates
Assuming the effective Hamiltonian (19) is valid, we can use it to implement 2-qubit quantum logic gates between two ions in a manner similar to the usual Cirac-Zoller (CZ) scheme . The analogy is not perfect because in the present case the Jaynes-Cummings Hamiltonian $`H_{2SS}`$ is valid only in the picture defined by the unitary operator in eq. (17), which varies according to which atom is being addressed. Before we realise a gate, we must first transform back into the ‘common’ picture (i.e., the one where the Hamiltonian in eq. (16) is defined) and see how the time evolution behaves there. In this case we have that an initial state $`|\psi \left(0\right)`$ evolves according to
$$|\psi \left(t\right)=V^{}\left(t\right)U_{JCM}\left(t\right)V\left(0\right)|\psi \left(0\right)$$
(22)
where $`V\left(t\right)`$ is given in eq. $`\left(\text{17}\right)`$ and $`U_{JCM}\left(t\right)=\mathrm{exp}\left(\frac{it}{\mathrm{}}H_{2SS}\right)`$. In particular, the following states have a simple time evolution:
$`||0`$ $``$ $`\mathrm{exp}\left({\displaystyle \frac{i\nu _qt}{2}}\right)||0`$ (23)
$`|+|0`$ $``$ $`e^{\frac{i\nu _qt}{2}}\mathrm{cos}\left({\displaystyle \frac{\nu _q\eta _{jq}t}{2}}\right)|+|0`$ (25)
$`e^{\frac{i\nu _qt}{2}}\mathrm{sin}\left({\displaystyle \frac{\nu _q\eta _{jq}t}{2}}\right)||1`$
$`||1`$ $``$ $`e^{\frac{i\nu _qt}{2}}\mathrm{cos}\left({\displaystyle \frac{\nu _q\eta _{jq}t}{2}}\right)||1+`$ (27)
$`+e^{\frac{i\nu _qt}{2}}\mathrm{sin}\left({\displaystyle \frac{\nu _q\eta _{jq}t}{2}}\right)|+|0`$
$`|+|1`$ $``$ $`e^{\frac{i\nu _qt}{2}}\mathrm{cos}\left({\displaystyle \frac{\nu _q\eta _{jq}t}{\sqrt{2}}}\right)|+|1`$ (29)
$`e^{\frac{i\nu _qt}{2}}\mathrm{sin}\left({\displaystyle \frac{\nu _q\eta _{jq}t}{\sqrt{2}}}\right)||2.`$
As we can see, we obtain the usual Jaynes-Cummings Rabi flops, except that here the atomic states for which the atom and mode dynamically entangle and disentangle themselves are the dressed states $`|\pm `$, not the bare states $`|g,|e`$. There are also some additional time-dependent phases.
In Appendix B, we demonstrate explicitly how this evolution can be used to implement a 2- qubit gate between two ions. We follow the same basic three-step pulse sequence proposed by Cirac and Zoller : first, a $`\pi `$-pulse is realised between ion 1 and the chosen vibrational ‘data bus’ mode, which is initially cooled to the ground state. This effectively maps the internal state onto the motional one and vice-versa, implementing the so-called SWAP gate. Second, a $`2\pi `$-pulse is applied between the mode and ion 2, realising an entangling gate between the two systems. Finally, a second $`\pi `$-pulse maps the motional state back onto the first ion, completing the ion-ion gate. In the lightshift scheme, some minor modifications in the sequence are necessary due to the fact that the ‘computational basis states’ of the ions (generally assumed to be the bare states $`|g,|e`$) are not favoured by the time-evolution above. This will then require a few extra $`1`$-qubit rotations in between the three basic steps. In the end, we are able to implement a C-NOT gate, with ion 2 acting as the ‘control’ qubit, using a sequence of six pulses (three 1-qubit and three 2-qubit pulses). In comparison, the original CZ proposal requires $`5`$ pulses to implement a C-NOT gate, with the ions assuming the opposite roles: ion 1 is the ‘control’ and ion 2 the ‘target’ qubit. We note that, in our protocol, some of the 1-qubit pulses may (at least in principle) be realised simultaneously with a 2-qubit pulse: pulses 1 and 2 in Appendix B can realised together, and the same is true of pulses 4 and 5. In contrast, in the CZ scheme each of the five pulses must be realised in sequence.
## III Comparative performance of gate schemes
We now study the performance of our ‘lightshift-based’ (LB) gate scheme, comparing it to that of Cirac and Zoller’s original ‘red-sideband pulse’ proposal . Briefly speaking, our goal is to estimate the overall switching rate for an ion-ion C-NOT gate that can likely be attained using each scheme.
We begin by recalling that this rate will be essentially governed by the speed of the three 2-qubit steps in either scheme’s pulse sequence. This follows since 1-qubit ionic gates are unlimited by the mode frequency, and can therefore be implemented at a much greater speed than 2-qubit ion-mode pulses . If we also assume for simplicity that the same ionic transition is used for both $`\pi `$ and $`2\pi `$pulses, then the overall ion-ion gate frequency should be approximately equal to the 2-qubit Jaynes-Cummings Rabi frequency. Here we are using the convention that one complete Rabi oscillation, i.e., when all states and their phases have returned to their initial values, corresponds to a $`4\pi `$ pulse.
For the LB scheme, this frequency is just $`\frac{\eta \nu }{2}`$. For a typical value $`\eta =0.1`$ of the Lamb-Dicke parameter, we obtain therefore an overall C-NOT switching rate of about $`\frac{\nu }{20}`$. Although still well under the limit posed by the mode frequency $`\nu `$ itself, such a rate would represent a substantial improvement with respect to current experiments. For example, in the (single-ion) 2-qubit gate experiment reported in , the 2-qubit Rabi frequency was approximately $`10^3\nu `$. In what follows, we elaborate on this comparison by making a more thorough analysis of the limits of validity of the two methods. In particular, we include numerical confirmation of the efficiency of the LB scheme.
### A Regimes of the Cirac-Zoller scheme
Unlike in the LB scheme, in the CZ method the speed of the 2-qubit gates is directly proportional to the laser field used to drive the red-sideband transition. This field cannot however be made too intense without driving unwanted off-resonant transitions, which therefore are the limiting factor on the resulting gate speed. Before we can properly assess this limit quantitatively, we must first recall that the CZ scheme operates in two strikingly different regimes, depending on the spatial profile of the laser field . The origin of this difference lies in the presence or not of strongly coupled off-resonant levels. It turns out that the conditions under which transitions to these levels can be safely ignored (as is implied in the derivation of the CZ scheme) depend crucially whether travelling-wave or standing-wave laser radiation is employed to drive the red-sideband transition
When a travelling beam is used, the closest-lying off-resonant transition is the carrier transition itself, which is detuned by the mode frequency $`\nu `$. Despite this, the carrier is also stronger than the resonant transition by a factor of $`\eta ^11`$. Intuitively, this situation is analogous to a V-type 3-level atom where a weak transition (of strength $`\eta \mathrm{\Omega }^{}`$) is being resonantly driven, and where there is another closely lying transition, detuned by $`\nu `$, which has a much stronger coupling constant $`\mathrm{\Omega }^{}`$ (Fig 3). The effects of both transitions must then be carefully weighed against each other: if $`\nu `$ is large with respect to $`\eta \mathrm{\Omega }^{}`$, then we may expect the off-resonant transition to be ‘washed out’ on average, as happens in usual rotating-wave approximations. However, this condition alone is not sufficient, since in the limit $`\eta 0`$ the off-resonant transition must dominate the time evolution, resulting in oscillations with effective Rabi frequency $`\frac{\mathrm{\Omega }^2}{\nu }`$. We can therefore expect the resonant transition to dominate only if its secular Rabi frequency $`\eta \mathrm{\Omega }^{}`$ is much greater than this value, i.e. if
$`\mathrm{\Omega }^{}\eta \nu .`$ (30)
The validity of this heuristic argument for the actual Cirac-Zoller Hamiltonian can be confirmed via a straightforward perturbation-theory calculation .
In other words, in order to ignore off-resonant transitions the Rabi frequency $`\mathrm{\Omega }^{}`$ of the ion-mode interaction must be extremely small, of the order $`\nu /100`$ for a typical value $`\eta =0.1`$. This in turn implies that the switching rate of the resulting logic gates will be of order $`\eta \mathrm{\Omega }^{}\nu /1000`$, way below the upper limit set by $`\nu `$. It is worthwhile to note that eq. (30) was indeed satisfied in both published experiments that implemented CZ-like Rabi flops using travelling-wave radiation and a single trapped ion .
A very different situation arises if the laser field forms a sinusoidal standing wave (such as could be obtained by bouncing the beam back on itself from a mirror), and if the ion is located exactly in one of the nodes of this wave. In this case, interference from the two travelling components of the wave completely cancels many of the off-resonant transitions, in particular the carrier . This effective selection rule greatly increases the laser power that can be used, since the most important off-resonant terms remaining in the Hamiltonian (Jaynes-Cummings counter-rotating terms and terms describing the accidental driving of the wrong modes) are no longer stronger than the resonant one. Standard perturbation-theoretic arguments show that in this case the laser power should satisfy
$`\mathrm{\Omega }^{}{\displaystyle \frac{\nu }{\eta }}.`$ (31)
For $`\eta =0.1`$, this implies an increase by two orders of magnitude with respect to the travelling-wave case. As a result, this configuration could potentially lend itself to the implementation of much faster gates than the ones already achieved experimentally. Unfortunately, the technical difficulty of reliably maintaining an ion precisely in a wave node seems to have discouraged researchers from attempting such an experiment . We are also not aware of any current plans for experiments in this direction.
### B Efficiency of gate implementations
In what follows, we compare our ‘lightshift-based’ proposal to both regimes of the CZ scheme. We find that its performance can approach that of the standing-wave CZ configuration, without the latter’s technical drawbacks. In other words, an improvement of over an order of magnitude in the switching rate can be achieved with respect to current travelling-wave-based experiments without a great change in the experimental setup itself. It must be emphasised again that we are only interested here in the theoretical limits to the gate performance, arising exclusively from the existence of stray off-resonant excitations in the system. In other words, we are not concerned with external noise or dissipative effects such as spontaneous emission , but with the maximum performance obtainable even under ideal experimental conditions.
A useful figure of merit for comparing the performance of the different schemes can be defined as follows. First, we determine how efficiently the $`\pi `$-pulse (or ‘SWAP gate’) step is implemented in each scheme as a function of a relevant external parameter of the system, for instance laser power (a precise definition of what me mean by ‘efficiency’ is given below). We can then define the maximum switching rate for each scheme as the greatest speed that can be attained while simultaneously keeping the efficiency above a sufficiently high threshold, which we (arbitrarily) set at 99%.
The definition of ‘efficiency’ is also somewhat arbitrary. We take it to be the average fidelity with which the SWAP gate operates, maximized over one cycle, or
$$F(\eta ,\mathrm{\Omega })=\mathrm{max}|_{\text{1st cycle}}\frac{1}{n}\underset{k=1}{\overset{n}{}}\left|\psi _f^k|U(\eta ,\mathrm{\Omega },t)|\psi _i^k\right|^2,$$
(32)
where the average is taken over some set of ‘relevant’ initial states $`\{\psi _i^k\}_{k=1}^n`$, with ideal images under SWAP given by $`\{\psi _f^k\}_{k=1}^n`$, and where $`U(\eta ,\mathrm{\Omega },t)`$ represents the full time evolution of the ion-trap system. For simplicity, we take this set to be the basis states $`\{|g|0,|g|1\}`$ (in the case of the CZ gate) or $`\{||0,||1\}`$ (in the case of the LB gate).
#### 1 Numerical results
In Fig. 4 we plot the efficiency function F($`\eta ,\mathrm{\Omega }^{}`$), with $`\eta `$ fixed at 0.1, for three different gate schemes: the CZ scheme using (a) travelling-wave radiation or (b) standing-wave radiation; and (c) the ‘lightshift-based’ scheme. The graphs were obtained by numerical integration of the full Schrödinger equation describing an ion-CM mode interaction in a two-ion trap, including all off-resonant transitions and all orders of the Lamb-Dicke parameter. The second or ‘stretch’ mode is assumed to be cooled to the ground state.
As should be expected, in the CZ schemes the efficiency decreases essentially monotonically with the laser power. In addition, the dramatic difference in performance between the standing- and travelling-wave CZ configurations is readily apparent (note the difference in scale of the two graphs). Indeed, if we consider 99% efficiency as the criterion for acceptable gate performance, then the upper limit for $`\mathrm{\Omega }^{}`$ in the travelling-wave case is about $`1.5\times 10^2\nu `$, while in the standing-wave case about $`1.25\nu `$, in agreement with the estimates in eqs. (30) and (31). Meanwhile, the efficiency of the LB scheme has a narrow peak around the resonance value $`\mathrm{\Omega }^{}=\frac{\nu }{2}`$, with a maximum value well over 0.99. (This is in good agreement with eq. (21), which predicts a population leakage of $`ϵ^20.005`$ into the stretch mode). The width of the region where $`F>0.99`$ is of the order of $`0.005\nu `$. We can conclude that highly efficient gate performance in this scheme is possible as long as the Rabi frequency of the laser-ion interaction is stable to within at least $`\pm 0.5\%`$.
### C Discussion
Our results indicate that, as long as the challenges of individual laser access and ground-state cooling can be met, the lightshift-based scheme should indeed allow highly efficient 2-qubit gates to be implemented within the Lamb-Dicke regime. Furthermore, the relatively high laser power employed in this scheme means these gates should be over an order of magnitude faster than their counterparts obtainable via the travelling-wave CZ scheme used in current experiments. Specifically, an ion-ion C-NOT gate with switching rate around $`\nu /20`$ may be realised. This speed is comparable to the one obtainable in principle with a standing-wave CZ configuration, but our proposal achieves it without requiring a precisely controlled standing-wave field. We believe that these features should make the lightshift-based scheme a attractive candidate for the realisation of faster quantum gates. Furthermore, testing the underlying principle of the scheme in existing single-ion traps should present no difficulty.
Finally, we would like to briefly compare our scheme with the ‘magic Lamb-Dicke parameter’ (MLDP) proposal of Monroe et al. . This elegant scheme exploits the fact that the 1-qubit Rabi frequency $`\mathrm{\Omega }`$ in eq. (3) is in fact dependent on the number of motional excitations of the ion. It turns out that, for specific ‘magic’ values of the Lamb-Dicke parameter, the values of $`\mathrm{\Omega }`$ corresponding to zero and one phonons become commensurate. This then means that, after a sufficient number of Rabi periods, the atomic state is flipped or not depending on the state of the mode, in other words a C-NOT gate with the mode as control qubit can be implemented. The scheme has a number of experimental advantages, notably the absence of the ‘auxiliary’ level needed in the CZ and LB schemes. Also, since it only uses the strong ionic ‘carrier’ transition, the laser power used can be quite considerable, leading also to relatively fast gates. The exact switching rate that can be obtained depends on the chosen ‘magic’ value, but should be as least as large as the ones obtained by the other methods discussed in this paper (see for a discussion).
The method however has also at least two drawbacks. First of all, even the smallest ‘magic’ value of $`\eta `$ quoted in is 0.316. This is already a bit too large for the validity of the Lamb-Dicke regime required by currently used cooling mechanisms such as sideband cooling . Unless more sophisticated cooling methods are employed (possibly involving the use of higher-order sidebands ), one would then need the ability to fine-tune $`\eta `$ to different values at different stages of the experiment, a feat that has not yet been accomplished in practice to our knowledge.
A second drawback comes the fact that the MLDP scheme can only implement universal ion-ion quantum logic if it is supplemented with another mechanism capable of realising SWAP gates between internal and motional states. For example, in , Monroe et al point out that an ion-ion C-NOT gate can be realising by “sandwiching” an MLDP-based ion-mode C-NOT between two SWAP gates, just as happens in the CZ scheme. However, the dispersive interaction exploited in the MLDP scheme does not itself allow the transfer of excitations from the internal to the motional states. This can be seen by noting that none of the available gates (1-qubit ionic rotations and C-NOT gates with the mode as control qubit) changes the populations in any motional state. (In other words, these operations alone do not constitute a universal set of gates . SWAP gates can only be realised via some different mechanism, for instance the CZ red-sideband method or our LB method. In particular, a gate using LB-based SWAP steps and an MLDP-based entangling step would combine the best features of both these schemes, including both speed and the absence of complications such as auxiliary levels and standing waves (Note though that, since the LB scheme requires a smaller value of $`\eta `$, the ability to tune this parameter would still be required). Whether in this “hybrid” combination or on its own, we hope that the LB scheme will prove to be a useful tool for ion-trap quantum information processing.
Endnote: Shortly after this work was submitted, another study of the speed limits of Cirac-Zoller gates was put forward by Steane et al. . Apart from presenting results which support and extend the discussion in section IIIA above, these authors also propose and experimentally test an independent method for increasing the gate switching rates within a travelling-wave scenario. Their idea is somewhat complementary to the one presented in this paper: they argue that the rapid decay in gate efficiency shown in fig. 4a) is partly due to a shift in the sideband transition frequency caused by the nearby strong carrier transition. This shift can be compensated for by choosing the laser beam to be slightly detuned from the first sideband frequency, resulting in gates that are considerably faster than the “standard” CZ gates we have considered in our analysis. Nevertheless it appears that, if a sufficiently high gate fidelity is demanded, then our lightshift-based scheme is still faster than even this enhanced scheme .
Acknowledgements: We thank Dana Berkeland, Chris Monroe, H. Christoph Nägerl, Juan Poyatos, Danny Segal, Andrew Steane, Jörg Steinbach, and Antônio Vidiella-Barranco for suggestions, clarifications and comments. We acknowledge the support of the Brazilian agency Conselho Nacional de Desenvolvimento Científico e Tecnológico (CNPQ), the ORS Award Scheme, the United Kingdom Engineering and Physical Sciences Research Council, the Leverhulme Trust, the European Science Foundation, and the European Union.
## A Limits to lightshift gates in N-ion strings
James has given detailed numerical data for the mode parameters of up to $`10`$ trapped ions. It turns out that the frequency $`\nu _q`$ of a mode of any given order $`q`$ is roughly independent of the number of ions (to about $`0.5\%`$ over the range of ion numbers investigated). In the second line of the table below we reproduce these rough frequency values for the first six modes, relative to the frequency $`\nu _1`$ of the lowest (CM) mode .
| $`q`$ | 1 | 2 | 3 | 4 | 5 | 6 |
| --- | --- | --- | --- | --- | --- | --- |
| $`\frac{\nu _q}{\nu _1}`$ | 1 | $`\sqrt{3}1.73`$ | 2.41 | 3.06 | 3.68 | 4.28 |
| $`\mathrm{min}|{}_{p}{}^{}\frac{\left|\nu _p\nu _q\right|}{\nu _q}`$ | $`0.73`$ | 0.39 | 0.27 | 0.20 | 0.16 | |
| $`\eta _{\mathrm{max}}`$ | 0.146 | 0.08 | 0.05 | 0.04 | 0.03 | |
| $`\frac{\eta _{\mathrm{max}}\nu _q}{2\nu _1}`$ | 0.073 | 0.069 | 0.065 | 0.061 | 0.055 | |
We can use this data along with the condition in eq. 21,i.e.
$$ϵ^2=\left(\frac{\eta _{jq}\nu _q}{2\left|\nu _p\nu _q\right|}\right)^21,$$
(A1)
in order to estimate the range of values of the Lamb-Dicke parameter $`\eta _{jq}`$ for which the lightshift-based scheme should work within a given precision. (It can be verified that losses due to other off-resonant transitions such as the counter-rotating terms in eq. (18) are relatively small in the limit of small $`\eta _{jq}`$). For each mode, we list in the third line the relative frequency spacing to its closest-lying neighbour. Note that the closest mode is always the next-highest one, and that their relative spacing decreases with increasing mode order. In the fourth line, we list the maximum value $`\eta _{\mathrm{max}}`$ that $`\eta _{jq}`$ can assume such that $`ϵ^20.01`$. Within this limit we should be able to discard all off-resonant terms in the Hamiltonian in eq. (18), and the dynamics is then well described by the effective Jaynes-Cummings interaction in eq. (19). Finally, in the fifth line we give the resulting maximum Rabi frequency achievable using each mode (relative to the CM mode frequency). Note that the increase of the mode frequencies themselves is completely compensated by the decrease in the allowed Lamb-Dicke parameters, with the effect that the overall Rabi frequency also diminishes as the mode order is increased.
## B C-NOT gate in the Lightshift scheme
The following sequence of pulses realises a C-NOT gate between the internal states of two trapped ions, using the LB ion-mode interaction given in eqs.(23-29).
1. First, assuming the ‘bus’ mode is initially in the ground state, the state of ion 1 in the $`|\pm _1`$ basis is mapped onto the $`|0`$ and $`|1`$ phonon states by a 2-qubit $`\pi `$-pulse of duration $`\tau _1=\frac{\pi }{\nu _q\eta _{qj}}`$
$`|_1|0\stackrel{\tau _1}{}e^{\frac{i\pi }{2\eta _{jq}}}|_1|0`$ (B2)
$`|+_1|0\stackrel{\tau _1}{}e^{\frac{i\pi }{2\eta _{jq}}}|_1|1.`$ (B3)
The phase is identical for both initial states and can be ignored; ion 1 is left in the $`|_1`$ state. In terms of the logical basis $`|g_1,|e_1`$, this transformation corresponds to applying a sequence of three gates: first a Hadamard rotation of the ion, followed by a SWAP gate with the mode, and finally a second Hadamard rotation.
2. A 1-qubit $`\frac{\pi }{2}`$ pulse coupling $`|g_2`$ to an unpopulated ‘auxiliary’ level $`|e^{}`$ is then applied on ion 2, mapping $`|g_2\frac{1}{\sqrt{2}}\left(|g_2|e^{}_2\right)|^{}_2`$. As in the CZ scheme, this $`|g_2|e^{}_2`$ transition should be chosen such that level $`|e_2`$ is not affected (for instance by using a different polarization).
3. A 2-qubit $`2\pi `$ pulse of duration $`\tau _2=\frac{2\pi }{\nu _q\eta _{jq}},`$ resonant with the $`|g_2|e^{}_2`$ transition, is applied on ion 2. States $`|e_2|0_2`$ and $`|e_2|1_2`$ of the ion-mode system are unaffected by this, while states $`|^{}_2|0_2,|^{}_2|1_2`$ evolve according to
$`|^{}_2|0_2\stackrel{\tau _2}{}\mathrm{exp}\left({\displaystyle \frac{i\pi }{\eta _{jq}}}\right)|^{}_2|0_2`$ (B5)
$`|^{}_2|1_2\stackrel{\tau _2}{}\mathrm{exp}\left({\displaystyle \frac{i\pi }{\eta _{jq}}}\right)|^{}_2|1_2.`$ (B6)
4. Another 1-qubit $`\frac{\pi }{2}`$ pulse coupling $`|g_2`$ to $`|e^{}_2`$ is then applied, mapping $`|^{}_2`$ back to $`|g_2.`$ For convenience, we assume here that this pulse also cancels the phase acquired in the previous step. The overall effect of the previous three pulses is to implement a ‘control-$`\sigma _z`$’ gate between the mode and ion2, which maps
$`|g_2|0`$ $``$ $`|g_2|0;|g_2|1|g_2|1`$ (B7)
$`|e_2|0`$ $``$ $`|e_2|0;|e_2|1|e_2|1.`$ (B8)
5. The state of the mode is then mapped back onto ion 1 by a second $`2`$-qubit $`\pi `$ pulse
$`|_1|0\stackrel{\tau _1}{}\mathrm{exp}\left({\displaystyle \frac{i\pi }{2\eta _{jq}}}\right)|_1|0`$ (B10)
$`|_1|1\stackrel{\tau _1}{}\mathrm{exp}\left({\displaystyle \frac{i\pi }{2\eta _{jq}}}\right)|+_1|0`$ (B11)
6. Finally, a 1-qubit pulse removes the phase acquired in the previous step, mapping states $`\mathrm{exp}\left({\displaystyle \frac{i\pi }{2\eta _{jq}}}\right)|\pm _1`$ of ion 1 into $`|_1`$. This completes the gate, whose overall effect in the computational basis is a C-NOT between ion 2 (the control qubit) and ion 1 (the target qubit):
$`|g_1|g_2|0\stackrel{1}{}|_1|g_2\left(|0|1\right)\stackrel{24}{}|_1|g_2\left(|0+|1\right)\stackrel{56}{}|g_1|g_2|0;`$ (13)
$`|g_1|e_2|0\stackrel{1}{}|_1|e_2\left(|0|1\right)\stackrel{24}{}|_1|e_2\left(|0|1\right)\stackrel{56}{}|e_1|e_2|0`$ (14)
$`|e_1|g_2|0\stackrel{1}{}|_1|g_2\left(|0+|1\right)\stackrel{24}{}|_1|g_2\left(|1|0\right)\stackrel{56}{}|e_1|g_2|0;`$ (15)
$`|e_1|e_2|0\stackrel{1}{}|_1|e_2\left(|0+|1\right)\stackrel{24}{}|_1|e_2\left(|0+|1\right)\stackrel{56}{}|g_1|e_2|0;`$ (16)
Note that the first five pulses already generate a ‘maximally entangling’ $`2`$ -qubit gate, equivalent to the C-NOT gate except for a local rotation. |
warning/0002/hep-th0002173.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Two dimensional integrable models represent an important laboratory for testing new ideas and developing new methods for constructing exact solutions as well as for the nonperturbative quantization of 4-D non-abelian gauge theories, gravity and string theory. Among the numerous technics for constructing 2-D integrable models and their solutions , , the two loop $`\widehat{𝒢}`$ -WZNW and gauged $`\widehat{G}/\widehat{H}`$-WZNW models have the advantage of providing a simple and universal method for derivation of the zero curvature representation as well as a consistent path integral formalism for their description. The power of such method was demonstrated in constructing (multi) soliton solution of the abelian affine Toda models and certain nonsingular nonabelian (NA) affine Toda models .
The present paper is devoted to the sistematic construction of the simplest class of singular torsionless affine NA Toda models characterized by the fact that the space of physical fields $`g_0^f`$ lies in the coset $`𝒢_0/𝒢_0^0=\frac{SL(2)}{U(1)}U(1)^{rank𝒢1}`$. Our main result is that such models exists only for the following three affine Kac-Moody algebras, $`B_n^{(1)},A_{2n}^{(2)}`$ and $`D_{n+1}^{(2)}`$ under certain specific restrictions on the choice of the subgroup $`𝒢_0^0`$ (i.e. equivalently the choice of gradation $`Q`$ and constant grade $`\pm 1`$ elements $`ϵ_\pm `$). It turns out that these models are T-selfdual (i.e. the axial and the vector gauging of the $`U(1)`$ factor in the coset $`\frac{SL(2)}{U(1)}U(1)^{rank𝒢1}`$ leads to identical actions ). They appear to be natural generalization of the Lund-Regge (“complex Sine-Gordon”) model and exactly reproduce the family of models proposed by Fateev . Our construction provide a simple proof of their classical integrability.
It is important to mention that relaxing the structure of the coset $`𝒢_0/𝒢_0^0=\frac{SL(2)U(1)^{rank𝒢1}}{U(1)}`$, i.e. gauging specific combinations of the Cartan subalgebra of $`𝒢_0`$ we find two new families of integrable models, axionic (for axial gauging) and torsionless ( for vector gauging) for all affine (twisted and untwisted ) Kac-Moody algebras which are T-dual (but not self dual).
An important motivation for the construction of the above singular NA Toda models is the fact that their soliton solutions (for imaginary coupling )carries both electric and magnetic (topological ) charges and have properties quite similar to the 4-D dyons of the Yang-Mills-Higgs model .
The generic NA Toda models are classified according to a $`𝒢_0𝒢`$ embedding induced by the grading operator $`Q`$ decomposing an finite or infinite Lie algebra $`𝒢=_i𝒢_i`$ where $`[Q,𝒢_i]=i𝒢_i`$ and $`[𝒢_i,𝒢_j]𝒢_{i+j}`$. A group element $`g`$ can then be written in terms of the Gauss decomposition as
$$g=NBM$$
(1.1)
where $`N=\mathrm{exp}𝒢_<H_{}`$, $`B=\mathrm{exp}𝒢_0`$ and $`M=\mathrm{exp}𝒢_>H_+`$. The physical fields $`B`$ lie in the zero grade subgroup $`𝒢_0`$ and the models we seek correspond to the coset $`H_{}\backslash G/H_+`$.
For consistency with the hamiltonian reduction formalism, the phase space of the G-invariant WZNW model is reduced by specifying the constant generators $`ϵ_\pm `$ of grade $`\pm 1`$. In order to derive an action for $`B𝒢_0`$, invariant under
$`gg^{}=\alpha _{}g\alpha _+,`$ (1.2)
where $`\alpha _\pm (z,\overline{z})H_\pm `$. we have to introduce a set of auxiliary gauge fields $`A𝒢_<`$ and $`\overline{A}𝒢_>`$ transforming as
$`AA^{}=\alpha _{}A\alpha _{}^1+\alpha _{}\alpha _{}^1,\overline{A}\overline{A}^{}=\alpha _+^1\overline{A}\alpha _++\overline{}\alpha _+^1\alpha _+.`$ (1.3)
The result is
$`S_{G/H}(g,A,\overline{A})`$ $`=`$ $`S_{WZNW}(g)`$
$``$ $`{\displaystyle \frac{k}{2\pi }}{\displaystyle 𝑑z^2Tr\left(A(\overline{}gg^1ϵ_+)+\overline{A}(g^1gϵ_{})+Ag\overline{A}g^1\right)}.`$
Since the action $`S_{G/H}`$ is $`H`$-invariant, we may choose $`\alpha _{}=N_{}^1`$ and $`\alpha _+=M_+^1`$. From the orthogonality of the graded subpaces, i.e. $`Tr𝒢_ij=0,i+j0`$, we find
$`S_{G/H}(g,A,\overline{A})`$ $`=`$ $`S_{G/H}(B,A^{},\overline{A}^{})`$ (1.4)
$`=`$ $`S_{WZNW}(B){\displaystyle \frac{k}{2\pi }}{\displaystyle 𝑑z^2Tr[A^{}ϵ_++\overline{A}^{}ϵ_{}+A^{}B\overline{A}^{}B^1]},`$
where
$`S_{WZNW}={\displaystyle \frac{k}{4\pi }}{\displaystyle d^2zTr(g^1gg^1\overline{}g)}{\displaystyle \frac{k}{24\pi }}{\displaystyle _D}ϵ_{ijk}Tr(g^1_igg^1_jgg^1_kg),`$ (1.5)
and the topological term denotes a surface integral over a ball $`D`$ identified as space-time.
Action (1.4) describe the non singular Toda models among which we find the Conformal and the Affine abelian Toda models where $`Q=_{i=1}^r\frac{2\lambda _iH}{\alpha _i^2},ϵ_\pm =_{i=1}^rc_{\pm i}E_{\pm \alpha _i}`$ and $`Q=hd+_{i=1}^r\frac{2\lambda _iH}{\alpha _i^2},ϵ_\pm =_{i=1}^rc_{\pm i}E_{\pm \alpha _i}^{(0)}+E_\psi ^{(\pm 1)}`$ respectively, where $`\psi `$ denote the highest root, $`\lambda _i`$, the fundamental weights, $`h`$ the coxeter number of $`𝒢`$ and $`H_i`$ are the Cartan subalgebra generators in the Cartan Weyl basis, $`Tr(H_iH_j)=\delta _{ij}`$.
More interesting cases arises in connection with non abelian embeddings $`𝒢_0𝒢`$. In particular, if we supress one of the fundamental weights from $`Q`$, the zero grade subspace acquires a nonabelian structure $`sl(2)u(1)^{rank𝒢1}`$. Let us consider for instance $`Q=h^{}d+_{ia}^r\frac{2\lambda _iH}{\alpha _i^2}`$, where $`h^{}=0`$ or $`h^{}0`$ corresponding to the Conformal and Affine nonabelian (NA) Toda respectively. The absence of $`\lambda _a`$ in $`Q`$ prevents the contribution of the simple root step operator $`E_{\alpha _a}^{(0)}`$ in constructing $`ϵ_+`$. It in fact, allows for reducing the phase space even further. This fact can be understood by enforcing the nonlocal constraint $`J_{YH}=\overline{J}_{YH}=0`$ where $`Y`$ is such that $`[YH,ϵ_\pm ]=0`$ and $`J=g^1g`$ and $`\overline{J}=\overline{}gg^1`$. Those generators of $`𝒢_0`$ commuting with $`ϵ_\pm `$ define a subalgebra $`𝒢_0^0𝒢_0`$. Such subsidiary constraint is incorporated into the action by requiring symmetry under
$`gg^{}=\alpha _0g\alpha _0^{}`$ (1.6)
where we shall consider $`\alpha _0^{}=\alpha _0(z,\overline{z})𝒢_0^0`$, i.e., axial symmetry (the vector gauging is obtained by choosing $`\alpha _0^{}=\alpha _{0}^{}{}_{}{}^{1}(z,\overline{z})𝒢_0^0`$). Auxiliary gauge fields $`A_0=a_0YH`$ and $`\overline{A}_0=\overline{a}_0YH𝒢_0^0`$ are introduced by splitting $`A=A_{}+A_0`$ and $`\overline{A}=\overline{A}_++\overline{A}_0`$ transforming as
$`AA^{}=\alpha _0A\alpha _0^1+\alpha _0\alpha _0^1,\overline{A}\overline{A}^{}=\alpha _0^1\overline{A}\alpha _0+\overline{}\alpha _0^1\alpha _0,`$
$`A_0A_0^{}=A_0+\alpha _0^1\alpha _0,\overline{A}_0\overline{A}_0^{}=\overline{A}_0+\overline{}\alpha _0\alpha _0^1.`$
The invariant action, under transformations generated by $`H_{(,0)}`$ (in left) and $`H_{(+,0)}`$ (in right) is given by
$`S_{G/H}(g,A,\overline{A})`$ $`=`$ $`S_{WZNW}(g){\displaystyle \frac{k}{2\pi }}{\displaystyle }d^2zTr(A(\overline{}gg^1ϵ_+)`$ (1.7)
$`+`$ $`\overline{A}(g^1gϵ_{})+Ag\overline{A}g^1+A_0\overline{A}_0).`$
Notice that the physical fields $`g_0^f`$ lie in the coset $`𝒢_0/𝒢_0^0=sl(2)u(1)^{rank𝒢1}/u(1)`$ of dimension $`rank𝒢+1`$ and are classified according to the gradation $`Q`$. It therefore follows that $`S_{G/H}(g,A,\overline{A})=S_{G/H}(g_0^f,A^{},\overline{A}^{})`$.
In a detailed study of the gauged WZNW construction for finite dimensional Lie algebras leading to Conformal NA Toda models was presented. The study of its symmetries was given in refs. and in . Here we generalize the construction of ref. to infinite dimensional Lie algebras leading to NA Affine Toda models characterized by the broken conformal symmetry and by the presence of solitons.
Consider the Kac-Moody algebra $`\widehat{𝒢}`$
$`[T_m^a,T_n^b]=f^{abc}T_{m+n}^c+\widehat{c}m\delta _{m+n}\delta ^{ab}`$
$`[\widehat{d},T_n^a]=nT_n^a;[\widehat{c},T_n^a]=[\widehat{c},\widehat{d}]=0`$ (1.8)
The NA Toda models we shall be constructing are associated to gradations of the type<sup>3</sup><sup>3</sup>3 We are considering $`𝒢^{}`$ to be semisimple and therefore $`a`$ to be one of the end points of the Dynkin diagram of $`𝒢`$, otherwise the model decomposes in two abelian Toda models coupled by $`\psi `$ and $`\chi `$.. $`Q_a(h^{})=h^{}d+_{ia}^r\frac{2\lambda _iH}{\alpha _i^2}`$, where $`h^{}`$ is chosen such that the zero grade subalgebra $`\widehat{𝒢}_0`$ defined by $`Q_a(h^{})`$, coincide with the zero grade subalgebra $`𝒢_0`$ defined by $`Q_a(h^{}=0)`$ (apart from two comuting generators $`\widehat{c}`$ and $`\widehat{d}`$). Since they commute with $`𝒢_0`$, the kinetic part decouples such that the conformal and the affine singular NA-Toda models differ only by the potential term. Due to the specific graded structure of the algebra, the following trace properties hold
$`\mathrm{Tr}A\overline{}g_0^f(g_0^f)^1=\mathrm{Tr}A_0\overline{}g_0^f(g_0^f)^1,\mathrm{Tr}\overline{A}(g_0^f)^1g_0^f=\mathrm{Tr}\overline{A}_0(g_0^f)^1g_0^f`$
$`\mathrm{Tr}Ag_0^f\overline{A}^{}(g_0^f)^1=\mathrm{Tr}A_0g_0^f\overline{A}_0(g_0^f)^1+\mathrm{Tr}A_{}g_0^f\overline{A}_+(g_0^f)^1.`$
Henceforth the action decomposes into three parts, i. e.,
$`S_{G/H}=S_{WZNW}(g_0^f)+F_0+F_\pm ,`$ (1.9)
where
$$F_0=\frac{k}{2\pi }d^2zTr\left(A_0\overline{}g_0^f(g_0^f)^1+\overline{A}_0(g_0^f)^1g_0^f+A_0g_0^f\overline{A}_0(g_0^f)^1+A_0\overline{A}_0\right),$$
$`F_\pm ={\displaystyle \frac{k}{2\pi }}{\displaystyle 𝑑z^2Tr\left(A_{}\widehat{ϵ_+}+\overline{A}_+\widehat{ϵ_{}}+A_{}g_0^f\overline{A}_+(g_0^f)^1\right)},`$
and the functional integral now factorizes into
$`Z={\displaystyle DA_0D\overline{A}_0\mathrm{exp}(F_0)DA_{}D\overline{A}_+\mathrm{exp}(F_\pm )}.`$ (1.10)
The integration over the auxiliary gauge fields $`A`$ and $`\overline{A}`$ require explicit parametrization of $`B`$.
$`B=\mathrm{exp}(\stackrel{~}{\chi }E_{\alpha _a})\mathrm{exp}(RY^jh_j+\mathrm{\Phi }(H)+\nu \widehat{c}+\eta \widehat{d})\mathrm{exp}(\stackrel{~}{\psi }E_{\alpha _a})`$ (1.11)
where $`\mathrm{\Phi }(H)=_{i=1}^{r1}\phi _iX_{}^{j}{}_{i}{}^{}h_j`$ , $`_{j=1}^rY^jX_{}^{j}{}_{i}{}^{}=0`$ and $`h_i=\frac{2\alpha _iH}{\alpha ^2},i=1,\mathrm{},r1`$. After gauging away the nonlocal field $`R`$, the factor group element becomes
$$g_0^f=\mathrm{exp}(\chi E_{\alpha _a})\mathrm{exp}(\mathrm{\Phi }(H)+\nu \widehat{c}+\eta \widehat{d})\mathrm{exp}(\psi E_{\alpha _a})$$
(1.12)
where $`\chi =\stackrel{~}{\chi }e^{\frac{1}{2}Y\alpha _a},\psi =\stackrel{~}{\psi }e^{\frac{1}{2}Y\alpha _a}`$ We therefore get for the zero grade component
$$F_0=\frac{k}{2\pi }\left(a_0\overline{a}_02Y^2\mathrm{\Delta }+2(\frac{\alpha _aY}{\alpha _a^2})(\overline{a}_0\psi \chi +a_0\chi \overline{}\psi )e^{\mathrm{\Phi }(\alpha _a)}\right)d^2x$$
(1.13)
where $`\mathrm{\Delta }=1+\frac{(Y\alpha _a)^2}{2Y^2}\psi \chi e^{\mathrm{\Phi }(\alpha _a)}`$ and
$$F_\pm =\frac{k}{2\pi }\left(Tr(A_{}g_0^f\widehat{ϵ_{}}(g_0^f)^1)g_0^f(\overline{A}_+(g_0^f)^1\widehat{ϵ_+}g_0^f)(g_0^f)^1\right)d^2x$$
(1.14)
The effective action is obtained by integrating over the auxiliary fields $`A_0,\overline{A}_0,A_{}`$ and $`\overline{A}_+`$,
$$Z_0=DA_0D\overline{A}_0\mathrm{exp}(F_0)e^{S_0}$$
(1.15)
where $`S_0=\frac{k}{\pi }(\frac{Y\alpha _a}{\alpha _a^2})\frac{\psi \chi \overline{}\psi \chi }{Y^2\mathrm{\Delta }}e^{2\mathrm{\Phi }(\alpha _a)}`$. Also,
$$Z_\pm =DA_{}D\overline{A}_+\mathrm{exp}(F_\pm )e^{\frac{k}{2\pi }{\scriptscriptstyle Tr\widehat{ϵ_+}g_0^f\widehat{ϵ_{}}(g_0^f)^1}}$$
(1.16)
The total action (1.7) is therefore given as
$$S=\frac{k}{4\pi }\left(Tr(\mathrm{\Phi }(H)\overline{}\mathrm{\Phi }(H))+\frac{2\overline{}\psi \chi }{\mathrm{\Delta }}e^{\mathrm{\Phi }(\alpha _a)}+\eta \overline{}\nu +\nu \overline{}\eta 2Tr(\widehat{ϵ_+}g_0^f\widehat{ϵ_{}}(g_0^f)^1)\right)$$
(1.17)
Note that the second term in (1.17) contains both symmetric and antisymmetric parts:
$`{\displaystyle \frac{e^{\mathrm{\Phi }(\alpha _a)}}{\mathrm{\Delta }}}\overline{}\psi \chi ={\displaystyle \frac{e^{\mathrm{\Phi }(\alpha _a)}}{\mathrm{\Delta }}}(g^{\mu \nu }_\mu \psi _\nu \chi +ϵ_{\mu \nu }_\mu \psi _\nu \chi ),`$ (1.18)
where $`g_{\mu \nu }`$ is the 2-D metric of signature $`g_{\mu \nu }=diag(1,1)`$, $`=_0+_1\overline{}=_0_1`$. For $`n=1`$ ($`𝒢A_1`$, $`\mathrm{\Phi }(\alpha _1)`$ is zero) the antisymmetric term is a total derivative:
$`ϵ_{\mu \nu }{\displaystyle \frac{_\mu \psi _\nu \chi }{1+\psi \chi }}={\displaystyle \frac{1}{2}}ϵ_{\mu \nu }_\mu \left(\mathrm{ln}\left\{1+\psi \chi \right\}_\nu \mathrm{ln}{\displaystyle \frac{\chi }{\psi }}\right),`$ (1.19)
and it can be neglected. This $`A_1`$-NA-Toda model (in the conformal case), is known to describe the 2-D black hole solution for (2-D) string theory . The $`G_n`$-NA conformal Toda model can be used in the description of specific (n+1)-dimensional black string theories , with n-1-flat and 2-non flat directions ($`g^{\mu \nu }G_{ab}(X)_\mu X^a_\nu X^b`$, $`X^a=(\psi ,\chi ,\phi _i)`$), containing axions ($`ϵ_{\mu \nu }B_{ab}(X)_\mu X^a_\nu X^b`$) and tachions ($`\mathrm{exp}\left\{k_{ij}\phi _j\right\}`$), as well.
It is clear that the presence of the $`e^{\mathrm{\Phi }(\alpha _a)}`$ in (1.17) is responsible for the antisymmetric tensor generating the axionic terms. On the other hand, notice that $`\mathrm{\Phi }(\alpha _a)`$ depend upon the subsidiary nonlocal constraint $`J_{YH}=\overline{J}_{YH}=0`$ and hence upon the choice of the vector Y. It is defined to be orthononal to all roots contained in $`ϵ_\pm `$.
In ref. , the most general constant grade one element $`ϵ_+`$ was analysed and the precise condition for the absence of axions, i.e. $`\mathrm{\Phi }(\alpha _a)=0`$, was established determining a subclass of torsionless conformal NA singular Toda models (no torsion theorem). For finite dimensional Lie algebras, it was shown in ref. that the absence of axions can only occur for $`G_n=B_n`$, $`a=n`$ and and $`ϵ_\pm =_{i=1}^{n2}c_{\pm i}E_{\pm \alpha _i}+d_\pm E_{\pm (\alpha _n+\alpha _{n1})}`$. In such case, $`𝒢_0^0`$ is generated by $`YH=(\frac{2\lambda _n}{\alpha _n^2}\frac{2\lambda _{n1}}{\alpha _{n1}^2})H`$ and $`\mathrm{\Phi }(H)=_{i=1}^{n2}\phi _ih_i+\phi _{}(\alpha _{n1}+\alpha _n)H`$. Due to the root structure of $`B_n`$, we verify that $`\mathrm{\Phi }(\alpha _n)=\alpha _n(\alpha _{n1}+\alpha _n)\phi _{}=0`$
In extending the no torsion theorem to infinite affine Lie algebras, $`h^{}`$ is chosen by defining the gradation $`Q_a(h^{})`$ such that preserves $`𝒢_0`$. We consider $`\widehat{ϵ}_+=ϵ_++E_\psi ^{(1)}`$ where $`\psi `$ is the highest root of $`𝒢`$. Since conformal and the affine models differ only by the potential term, the solution for the no torsion condition is also satisfied for infinite dimensional algebras, whose Dynkin diagram possess a $`B_n`$-“tail like”. An obvious solution is the untwisted $`B_n^{(1)}`$ model. Two other solutions were found within the twisted affine Kac-Moody algebras $`A_{2n}^{(2)}`$ and $`D_{n+1}^{(2)}`$ as we shall describe in detail.
## 2 The $`B_n^{(1)}`$ Torsionless NA Toda model
Let $`Q=2(n1)d+_{i=1}^{n1}\frac{2\lambda _iH}{\alpha _i^2}`$ decomposing $`B_n^{(1)}`$ into graded subspaces. In particular $`𝒢_0=SL(2)U(1)^{n1}U(1)_{\widehat{c}}U(1)_{\widehat{d}}`$ generated by $`\{E_{\pm \alpha _n}^{(0)},h_1,\mathrm{},h_n,\widehat{c},\widehat{d}\}`$. Following the no torsion theorem of ref. , we have to choose $`\widehat{ϵ_\pm }=_{i=1}^{n2}c_{\pm i}E_{\pm \alpha _i}^{(0)}+c_{\pm (n1)}E_{\pm (\alpha _{n1}+\alpha _n)}^{(0)}+c_{\pm n}E_\psi ^{(\pm 1)}`$, where $`\psi =\alpha _1+2(\alpha _2+\mathrm{}+\alpha _{n1}+\alpha _n)`$ is the highest root of $`B_n`$ and $`𝒢_0^0`$ is generated by $`YH=(\frac{2\lambda _n}{\alpha _n^2}\frac{2\lambda _{n1}}{\alpha _{n1}^2})H`$ such that $`[𝐘𝐇,\widehat{ϵ_\pm }]=0`$. The coset $`𝒢_0/𝒢_0^0`$ is then parametrized according to (1.12) with $`\mathrm{\Phi }(H)=_{i=1}^{n1}_i\phi _i+\eta \widehat{h}+\nu \widehat{c}`$ where $`_i=(\alpha _n+\mathrm{}\alpha _i)H`$ so that $`Tr(_i_j)=\delta _{ij},i,j=1,\mathrm{},n1`$ and the total effective action becomes
$$S=\frac{k}{4\pi }d^2x\left(\frac{1}{4}\underset{i=1}{\overset{n1}{}}g^{\mu \nu }_\mu \phi _i_\nu \phi _i+g^{\mu \nu }\frac{_\mu \psi _\nu \chi }{1+\psi \chi }+\frac{1}{2}g^{\mu \nu }_\mu \nu _\nu \eta 2V\right)$$
(2.20)
where the “affine potential” $`(n>2)`$ is
$$V=\underset{i=1}{\overset{n2}{}}|c_i|^2e^{\phi _i\phi _{i+1}}+|c_{n1}|^22(1+2\psi \chi )e^{\phi _{n1}}+|c_n|^2e^{\phi _1+\phi _2\eta }$$
(2.21)
The action (2.20) is invariant under conformal tranformation
$`zf(z),\overline{z}g(\overline{z}),\psi \psi ,\chi \chi ,`$
$`\phi _s\phi _s+s\mathrm{ln}f^{}g^{};s=1,2,\mathrm{},n1;\eta \eta +2(n1)\mathrm{ln}f^{}g^{}`$ (2.22)
We should point out that the $`\eta `$ field plays a crucial role in establishing the conformal invariance of the theory. Integrable deformation of such class of theories can than, be sistematically obtained by setting $`\eta =0`$.
For the case $`n=2`$ we choose, $`\widehat{ϵ_\pm }=E_{\alpha _1+\alpha _2}^{(0)}+E_{\alpha _1\alpha _2}^{(1)}`$, $`\mathrm{\Phi }(\alpha _{n1})=\phi `$, i.e. $`\widehat{𝒢}=\widehat{S}O(5)`$, is also special in the sense that its complexified theory, i.e.
$`\psi i\psi ;\chi i\psi ^{};\phi i\phi `$
leads to the real action
$$S=\frac{k}{4\pi }d^2x\left(\frac{g^{\mu \nu }_\mu \psi _\nu \psi ^{}}{(1\psi \psi ^{})}+\frac{1}{4}g^{\mu \nu }_\mu \phi _\nu \phi +8(12\psi \psi ^{})cos\phi \right)$$
(2.23)
## 3 The twisted NA Toda Models
The twisted affine Kac-Moody algebras are constructed from a finite dimensional algebra possessing a nontrivial symmetry of their Dynkin diagrams (folding). Such symmetry can be extended to the algebra by an outer automorphism $`\sigma `$ , as
$$\sigma (E_\alpha )=\eta _\alpha E_{\sigma (\alpha )}$$
(3.24)
where $`\eta _\alpha =\pm 1`$. For the simple roots, $`\eta _{\alpha _i}=1`$. The signs can be consistently assign to all generators since nonsimple roots can be written as sum of two roots other roots.
The no torsion theorem require a $`B_n`$-“tail like” structure which is fulfilled only by the $`A_{2n}^{(2)}`$ and $`D_{n+1}^{(2)}`$ (see appendix N of ref. ). In both cases the automorphism is of order 2 (i.e. $`\sigma ^2=1`$).
Let us denote by $`\alpha `$ the roots of the untwisted algebra $`𝒢`$. For the $`A_{2n}^{(2)}`$ case, the automorphism is defined by
$$\sigma (\alpha _1)=\alpha _{2n},\sigma (\alpha _2)=\alpha _{2n1}\mathrm{},\sigma (\alpha _{n1})=\alpha _n$$
(3.25)
whilst for the $`D_{n+1}^{(2)}`$, the automorphism acts only in the “fish tail” of the Dynkin diagram of $`D_{n+1}`$, i.e.
$$\sigma (E_{\alpha _1})=E_{\alpha _1},\mathrm{},\sigma (E_{\alpha _{n1}})=E_{\alpha _{n1}},\sigma (E_{\alpha _n})=E_{\alpha _{n+1}}$$
(3.26)
The automorphism $`\sigma `$ decomposes the algebra $`𝒢=𝒢_{even}𝒢_{odd}`$. The twisted affine algebra is constructed from $`𝒢`$ assigning an affine index $`mZ`$ to the generators in $`𝒢_{even}`$ while $`mZ+\frac{1}{2}`$ to those in $`𝒢_{odd}`$ (see appendix N of ).
The simple root step operators for $`A_{2n}^{(2)}`$ are
$$E_{\beta _i}=E_{\alpha _i}^{(0)}+E_{\alpha _{2ni+1}}^{(0)},i=1,\mathrm{},nE_{\beta _0}=E_{\alpha _1\mathrm{}\alpha _{2n}}^{(\frac{1}{2})}$$
(3.27)
corresponding to the simple and highest roots
$$\beta _i=\frac{1}{2}(\alpha _i+\alpha _{2ni+1})i=1,\mathrm{},n,\psi =\alpha _1+\mathrm{}+\alpha _{2n}=2(\beta _1+\mathrm{}\beta _n)$$
(3.28)
respectively.
For $`D_{n+1}^{(2)}`$, simple root step operators are
$`E_{\beta _i}=E_{\alpha _i}^{(0)},i=1,\mathrm{},n1,E_{\beta _n}=E_{\alpha _n}^{(0)}+E_{\alpha _{n+1}}^{(0)}`$
$`E_{\beta _0}=E_{\alpha _1\mathrm{}\alpha _{n1}\alpha _{n+1}}^{(\frac{1}{2})}E_{\alpha _1\mathrm{}\alpha _{n1}\alpha _n}^{(\frac{1}{2})}`$ (3.29)
corresponding to the simple and highest roots
$`\beta _i=\alpha _ii=1,\mathrm{},n1,\beta _n={\displaystyle \frac{1}{2}}(\alpha _n+\alpha _{n+1}),`$
$$\psi =\alpha _1+\mathrm{}\alpha _{n1}+\frac{1}{2}(\alpha _n+\alpha _{n+1})=\beta _1+\mathrm{}\beta _n$$
(3.30)
where have denoted by $`\beta `$ the roots of the twisted (folded) algebra.
The torsionless affine NA Toda models are defined by
$$Q=2(2n1)\widehat{d}+\underset{in,n+1}{\overset{2n}{}}\frac{2\lambda _iH}{\alpha _i^2},$$
(3.31)
and
$$Q=(2n2)\widehat{d}+\underset{i=1}{\overset{n}{}}\frac{2\lambda _iH}{\alpha _i^2}$$
(3.32)
for $`A_{2n}^{(2)}`$ and $`D_{n+1}^{(2)}`$ respectively, where $`\lambda _i`$ are the fundamental weights of the untwisted algebra $`𝒢`$, i.e. $`\frac{2\lambda _i\alpha _j}{\alpha _j^2}=\delta _{ij}`$.
Both models are specified by the constant grade $`\pm 1`$ operators $`ϵ_\pm `$
$$\widehat{ϵ_\pm }=\underset{i=1}{\overset{n2}{}}c_{\pm i}E_{\pm \beta _i}+c_{\pm (n1)}E_{\pm (\beta _{n1}+\beta _n)}+c_{\pm n}E_{\beta _0}$$
(3.33)
where $`\beta _i`$ are the simple roots of the twisted affine algebra specified in (3.28) and in (3.30).
According to the grading generators (3.31) and (3.32), the zero grade subalgebra is in both cases $`𝒢_0=SL(2)U(1)^{n1}U(1)_{\widehat{c}}U(1)_{\widehat{d}}`$ generated by $`\{E_{\pm \beta _n}^{(0)},h_1,\mathrm{},h_n,\widehat{c},\widehat{d}\}`$. Hence the zero grade subgroup is parametrized as in (1.11) where we have taken $`\eta =0`$, responsible for breaking the conformal invariance. The factor group is given in (1.12), where $`𝒢_0^0`$ is generated by $`YH=(\frac{2\mu _n}{\beta _n^2}\frac{2\mu _{n1}}{\beta _{n1}^2})H`$ and $`\mu _i`$ are the fundamental weights of the twisted algebra i.e. $`\frac{2\mu _i\beta _j}{\beta _j^2}=\delta _{ij}`$ In order to decouple the $`\phi _i,i=1,\mathrm{},n1`$ we chose an orthonormal basis for the Cartan subalgebra, i.e. $`\mathrm{\Phi }(H)=_i\phi _i+\eta \widehat{h}+\nu \widehat{c}`$ where
$$_i=(\alpha _i+\mathrm{}\alpha _{2ni+1})H,YH=_n,Tr(_i_j)=2\delta _{ij},i,j=1,\mathrm{}n$$
(3.34)
and
$$_i=(\alpha _{ni+1}+\mathrm{}+\alpha _{n+1})H,YH=_n,Tr(_i_j)=\delta _{ij},i,j=1,\mathrm{}n$$
(3.35)
for $`A_{2n}^{(2)}`$ and $`D_{n+1}^{(2)}`$ respectively.
The Lagrangean density is obtained from (1.17) leading to
$$_{A_{2n}^{(2)}}=\frac{\chi \overline{}\psi }{1+\frac{1}{2}\psi \chi }+\frac{1}{2}\underset{i=1}{\overset{n1}{}}\phi _i\overline{}\phi _iV_{A_{2n}^{(2)}}$$
(3.36)
and
$$_{D_{n+1}^{(2)}}=2\frac{\chi \overline{}\psi }{1+\psi \chi }+\frac{1}{2}\underset{i=1}{\overset{n1}{}}\phi _i\overline{}\phi _iV_{D_{n+1}^{(2)}}$$
(3.37)
where
$$V_{A_{2n}^{(2)}}=\underset{i=1}{\overset{n2}{}}|c_i|^2e^{\phi _i+\phi _{i+1}}+\frac{1}{2}|c_n|^2e^{2\phi _1}+|c_{n1}|^2e^{\phi _{n1}}(1+\psi \chi )$$
(3.38)
and
$$V_{D_{n+1}^{(2)}}=\underset{i=1}{\overset{n2}{}}|c_i|^2e^{\phi _i+\phi _{i+1}}+\frac{1}{2}|c_n|^2e^{\phi _1}+|c_{n1}|^2e^{\phi _{n1}}(1+2\psi \chi )$$
(3.39)
The models described by (2.20), (3.36) and (3.37) coincide with those proposed by Fateev in .
## 4 Zero Curvature
The equations of motion for the NA Toda models are known to be of the form
$$\overline{}(B^1B)+[\widehat{ϵ_{}},B^1\widehat{ϵ_+}B]=0,(\overline{}BB^1)[\widehat{ϵ_+},B\widehat{ϵ_{}}B^1]=0$$
(4.40)
The subsidiary constraint $`J_{YH}=Tr(B^1BYH)=\overline{J}_{YH}=Tr(\overline{}BB^1YH)=0`$ can be consistenly imposed since $`[YH,\widehat{ϵ_\pm }]=0`$ as can be obtained from (4.40) by taking the trace with $`Y.H`$. Solving those equations for the nonlocal field $`R`$ yields,
$$R=(\frac{Y\alpha _n}{Y^2})\frac{\psi \chi }{\mathrm{\Delta }}e^{\mathrm{\Phi }(\alpha _n)},\overline{}R=(\frac{Y\alpha _n}{Y^2})\frac{\chi \psi }{\mathrm{\Delta }}e^{\mathrm{\Phi }(\alpha _n)}$$
(4.41)
The equations of motion for the fields $`\psi ,\chi `$ and $`\phi _i,i=1,\mathrm{},n1`$ obtained from (4.40) imposing the constraints (4.41) coincide precisely with the Euler-Lagrange equations derived from (3.36) and (3.37). Alternatively, (4.40) admits a zero curvature representation $`\overline{A}\overline{}A+[A,\overline{A}]=0`$ where
$$A=\widehat{ϵ_{}}+B^1B,\overline{A}=B^1\widehat{ϵ_+}B$$
(4.42)
Whenever the constraints (4.41) are incorporated into $`A`$ and $`\overline{A}`$ in (4.42), equations (4.40) yields the zero curvature representation of the NA singular Toda models. Such argument is valid for all NA Toda models, in particular for the torsionless class of models discussed in the previous two sections.
Using the explicit parametrization of $`B`$ given in (1.11), the corresponding $`\widehat{ϵ_\pm }`$ specified in (3.33), (3.28) and (3.30) together with (4.41) where $`Y`$ is given in (3.34) and (3.35), we find, in a systematic manner, the following form for $`A`$ and $`\overline{A}`$
$`A_{A_{2n}^{(2)}}`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{n2}{}}}c_i(E_{\alpha _i}^{(0)}+E_{\alpha _{2ni+1}}^{(0)})+c_{n1}(E_{\alpha _n\alpha _{n1}}^{(0)}+E_{\alpha _{n+1}\alpha _{n+2}}^{(0)})`$ (4.43)
$`+`$ $`c_nE_{\alpha _1+\mathrm{}+\alpha _{2n}}^{(\frac{1}{2})}+\psi e^{\frac{1}{2}R}(E_{\alpha _n}^{(0)}+E_{\alpha _{n+1}}^{(0)})+{\displaystyle \underset{i=1}{\overset{n1}{}}}\phi _i_i`$
$`+`$ $`{\displaystyle \frac{\chi }{\mathrm{\Delta }}}e^{\frac{1}{2}R}(E_{\alpha _n}^{(0)}+E_{\alpha _{n+1}}^{(0)})`$
and
$`\overline{A}_{A_{2n}^{(2)}}`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{n2}{}}}c_ie^{\phi _i+\phi _{i+1}}(E_{\alpha _i}^{(0)}+E_{\alpha _{2ni+1}}^{(0)})+c_ne^{2\phi _1}E_{\alpha _1\mathrm{}\alpha _{2n}}^{(\frac{1}{2})}`$ (4.44)
$`+`$ $`c_{n1}e^{\phi _{n1}}(E_{\alpha _n+\alpha _{n1}}^{(0)}+E_{\alpha _{n+1}+\alpha _{n+2}}^{(0)})`$
$`+`$ $`c_{n1}\psi e^{\frac{1}{2}R\phi _{n1}}(E_{\alpha _{n+1}+\alpha _n+\alpha _{n1}}^{(0)}E_{\alpha _n+\alpha _{n+1}+\alpha _{n+2}}^{(0)})`$
$`+`$ $`c_{n1}\chi e^{\frac{1}{2}R\phi _{n1}}(E_{\alpha _{n1}}^{(0)}E_{\alpha _{n+2}}^{(0)})+c_{n1}\psi \chi e^{\phi _{n1}}(E_{\alpha _n+\alpha _{n1}}^{(0)}E_{\alpha _{n+1}+\alpha _{n+2}}^{(0)})`$
$`+`$ $`{\displaystyle \frac{1}{2}}c_{n1}\psi ^2\chi e^{\phi _{n1}\frac{1}{2}R}(E_{\alpha _{n+1}+\alpha _n+\alpha _{n1}}^{(0)}E_{\alpha _n+\alpha _{n+1}+\alpha _{n+2}}^{(0)})`$
$`A_{D_{n+1}^{(2)}}`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{n2}{}}}c_iE_{\alpha _i}^{(0)}+c_{n1}(E_{\alpha _n\alpha _{n1}}^{(0)}+E_{\alpha _{n1}\alpha _{n+1}}^{(0)})`$ (4.45)
$`+`$ $`c_n(E_{(\alpha _1+\mathrm{}+\alpha _{n1}+\alpha _{n+1})}^{(\frac{1}{2})}E_{(\alpha _1+\mathrm{}+\alpha _n+\alpha _{n+1})}^{(\frac{1}{2})})+\psi e^{\frac{1}{2}R}(E_{\alpha _n}^{(0)}+E_{\alpha _{n+1}}^{(0)})`$
$`+`$ $`{\displaystyle \underset{i=1}{\overset{n1}{}}}\phi _i_i+{\displaystyle \frac{\chi }{\mathrm{\Delta }}}e^{\frac{1}{2}R}(E_{\alpha _n}^{(0)}+E_{\alpha _{n+1}}^{(0)})`$
and
$`\overline{A}_{D_{n+1}^{(2)}}`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{n2}{}}}c_ie^{\phi _i+\phi _{i+1}}E_{\alpha _i}^{(0)}+c_{n1}e^{\phi _{n1}}(E_{\alpha _n+\alpha _{n1}}^{(0)}+E_{\alpha _{n+1}+\alpha _{n1}}^{(0)})`$ (4.46)
$`+`$ $`2c_{n1}\psi e^{\frac{1}{2}R\phi _{n1}}E_{\alpha _{n+1}+\alpha _n+\alpha _{n1}}^{(0)}+2c_{n1}\chi e^{\frac{1}{2}R\phi _{n1}}E_{\alpha _{n1}}^{(0)}`$
$`+`$ $`2c_{n1}\psi \chi e^{\phi _{n1}}(E_{\alpha _{n+1}+\alpha _{n1}}^{(0)}+E_{\alpha _{n1}+\alpha _n}^{(0)})+c_{n1}\psi ^2\chi e^{\frac{1}{2}R\phi _{n1}}E_{\alpha _{n+1}+\alpha _n+\alpha _{n1}}^{(0)}`$
$`+`$ $`c_{n+1}e^{\phi _1}(E_{(\alpha _1+\mathrm{}+\alpha _{n1}+\alpha _{n+1})}^{(\frac{1}{2})}E_{(\alpha _1+\mathrm{}+\alpha _n+\alpha _{n+1})}^{(\frac{1}{2})})`$
For the untwisted affine $`B_n^{(1)}`$ model of the previous section the zero curvature representation is obtained from
$`A_{B_n^{(1)}}`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{n2}{}}}c_iE_{\alpha _i}^{(0)}+c_{n1}E_{\alpha _n\alpha _{n1}}^{(0)}+c_nE_{\alpha _1+2(\alpha _2+\mathrm{}+\alpha _n)}^{(1)}`$ (4.47)
$`+`$ $`\psi e^{\frac{1}{2}R}E_{\alpha _n}^{(0)}+{\displaystyle \underset{i=1}{\overset{n1}{}}}\phi _i_i+{\displaystyle \frac{\chi }{\mathrm{\Delta }}}e^{\frac{1}{2}R}E_{\alpha _n}^{(0)}`$
$`\overline{A}_{B_n^{(1)}}`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{n2}{}}}c_ie^{\phi _i+\phi _{i+1}}E_{\alpha _i}^{(0)}+c_ne^{\phi _1+\phi _2}E_{(\alpha _1+2(\alpha _2+\mathrm{}+\alpha _n))}^{(1)}+2\chi e^{\phi _{n1}+\frac{1}{2}R}E_{\alpha _{n1}}^{(0)}`$ (4.48)
$`+`$ $`c_{n1}(1+2\psi \chi )e^{\phi _{n1}}E_{\alpha _{n1}+\alpha _n}^{(0)}2c_{n1}e^{\phi _{n1}\frac{1}{2}R}\psi (1+\psi \chi )E_{\alpha _{n1}+2\alpha _n}^{(0)}`$
The zero curvature representation of such subclass of torsionless NA Toda models shows that they are in fact classically integrable field theories. The construction of the previous sections provides a sistematic affine Lie algebraic structure underlying those models.
## 5 Conclusions
We have constructed a class of affine NA Toda models from the gauged two-loop WZNW models in which left and right symmetries are incorporated by a suitable choice of grading operator $`Q`$. Such framework is specified by grade $`\pm 1`$ constant generators $`ϵ_\pm `$ and the pair $`(Q,ϵ_\pm )`$ determines the model in terms of a zero grade subgroup $`𝒢_0`$. We have shown that for non abelian $`𝒢_0`$, it is possible to reduce even further the phase space by constraining to zero the currents commuting with $`ϵ_\pm `$ , ($`J𝒢_0^0)`$) to the fields lying in the coset $`𝒢_0/𝒢_0^0`$ only. Moreover, we have found a Lie algebraic condition which defines a class of T-selfdual torsionless models, for the case $`𝒢_0^0=U(1)`$. The action for those models were sistematicaly constructed and shown to coincide with the models proposed by Fateev , describing the strong coupling limit of specific 2-d models representing sine-Gordon interacting with Toda-like models. Their weak coupling limit appears to be the Thirring model coupled to certain affine Toda theories .
Following the same line of arguments of the previous sections, one can construct more general models, say, $`𝒢_0/𝒢_0^0=\frac{SL(2)U(1)^{n1}}{U(1)^s}`$, $`𝒢_0/𝒢_0^0=\frac{SL(2)SL(2)U(1)^{n2}}{U(1)}`$, $`𝒢_0/𝒢_0^0=\frac{SL(3)U(1)^{n2}}{U(1)}`$, etc. Those models represent more general NA affine Toda models obtained by considering specific gradations $`Q_{a,b,\mathrm{}}=h_{a,b,\mathrm{}}d+_{ia,b\mathrm{}}^n\frac{2\lambda _iH}{\alpha _i^2}`$. However the important problem of the classification of all integrable models obtained as gauged two loop $`G`$-WZNW models remains open.
Acknowledgments We are grateful to FAPESP and CNPq for financial support. |
warning/0002/hep-ph0002011.html | ar5iv | text | # Bottom quark electroproduction in variable flavor number schemes
## I Introduction
The H1 and ZEUS experiments at HERA now have enough integrated luminosity to study bottom quark deep inelastic electroproduction. Therefore it is an appropriate time to present up-to-date predictions for the bottom quark components of the deep inelastic structure functions $`F_i(x,Q^2)`$ where $`i=2,L`$. We assume that the bottom quark is produced in an extrinsic fashion and that the neutral current reaction dominates over the charged current one. This means that in fixed order perturbative QCD the heavy quark structure functions $`F_{i,b}(x,Q^2,m_b^2)`$, $`i=2,L`$ are given by the virtual-photon gluon-fusion processes and their higher order corrections with only light partons in the initial state. Notice that in the case of bottom quark production the light partons are represented by the gluon and the four light flavors u, d, s, c together with their anti-particles.
In the literature one has adopted two different treatments of extrinsic bottom quark production, which are known as the massive and massless descriptions. The former treats the bottom quark as a heavy quark (with mass $`m_b`$) and the partonic cross sections (or heavy quark coefficient functions) are described by fixed order perturbation theory (FOPT) as mentioned above. Notice that due to the work in the perturbation series is now known up to second order. The latter treatment, which has been rather popular among groups which fit parton densities to experimental data, treats the bottom quark as a massless quark so that it can be represented by a scale dependent parton density $`f_b(x,\mu ^2)`$. Although at first sight these approaches are completely different they are actually intimately related. It was shown in that the large logarithms of the type $`\mathrm{ln}(Q^2/m_b^2)`$, which appear in FOPT when $`Q^2m_b^2`$, can be resummed in all orders. The upshot of this procedure is that the bottom components of the deep inelastic structure functions $`F_{i,b}(x,Q^2,m_b^2)`$, where $`i=2,L`$, which in the FOPT approach are written as convolutions of heavy quark coefficient functions with four-flavor light-mass (u,d,s,c) parton densities, become, after resummation, convolutions of light-mass parton coefficient functions with five-flavor light-mass parton densities which also include a bottom quark density. This procedure leads to the so-called zero mass variable flavor number scheme (ZM-VFNS) for $`F_{i,b}(x,Q^2)`$ where the mass of the bottom quark is absorbed into the new five flavor densities. To implement this scheme one has to be careful to use quantities which are collinearly finite in the limit $`m_b0`$. From the above considerations it is clear that the FOPT approach is better when the bottom quark pair is produced near threshold (where $`Q^2(x^11)4m_b^2`$) because terms in $`m_b`$ are important in this kinematic region. On the other hand far above threshold, where also $`Q^2m_b^2`$, the large logarithms mentioned above dominate the structure functions so that the ZM-VFNS approach should be more appropriate. Both approaches are characterized by the number of active flavors involved in the description of the parton densities which are given by four and five respectively. Each scheme has different densities but the momentum sum rule either gets contributions from four-flavor densities or five-flavor densities and is always satisfied.
As most of the experimental data will be in the kinematical regime which is between the threshold and the large $`Q^2`$ region, a third approach, called the variable flavour number scheme (VFNS), has been introduced to describe the heavy quark components $`F_{i,b}(x,Q^2)`$ of the deep inelastic structure functions. Actually there are several such schemes. They include the Aivasis, Collins, Olness, Tung (ACOT) , scheme, the Buza, Matiounine, Smith, van Neerven (BMSN) , scheme, the Thorne, Roberts (TR) scheme and the Chuvakin, Smith, van Neerven (CSN) scheme. A discussion of them is given in the last reference. The difference between the schemes can be attributed to two ingredients entering in their construction. The first one is the mass factorization procedure carried out before the large logarithms can be resummed. The second one is the matching condition imposed on the heavy quark density, which has to vanish in the threshold region of the production process. Another aspect of these approaches is that one needs two sets of parton densities. For bottom quark production one set only contains densities in a four-flavor number scheme whereas the second one, which also includes a bottom quark density, is parametrized in a five-flavor number scheme. Both parameterizations have to satisfy the matching relations quoted in . Up to next-to-leading order (NLO) they are continuous at the scale $`\mu =m_b`$ whereas in next-to-next-to leading order (NNLO) the parton densities become discontinuous while going from a four to a five flavor scheme. Starting from a three-flavor set of parton densities given in we have recently constructed in a four-flavor set of densities which satisfied the matching relations in at the scale $`\mu =m_c`$. Then we evolved these densities with LO or NLO splitting functions up to the scale $`\mu =m_b`$ and constructed a five-flavor set which also satisfied the matching relations in . This set was further evolved with LO and NLO splitting functions up to high scales. Notice that since the NNLO splitting functions are unknown the only difference between the NLO and NNLO parton densities can be attributed to the boundary conditions at $`\mu =m_c`$ and $`\mu =m_b`$ where the latter densities become discontinuous contrary to the LO and NLO ones. We can now use these densities to discuss VFNS for bottom quark deep inelastic electroproduction, in particular in the CSN and BMSN schemes. The previous discussions in were focussed on applications to charm quark electroproduction.
Since any description for bottom quarks follows closely that for charm quarks we refer the interested reader to for most of the details and simply specialize to bottom quark electroproduction in Sec.II. We work to second order in the running coupling constant $`\alpha _s(\mu ^2)`$. Numerical results are shown for the structure functions $`F_{2,b}`$ and $`F_{L,b}`$ in the CSN and BMSN schemes.
## II Bottom quark structure functions
In this Section we consider bottom quark deep inelastic electroproduction in two variable flavor number schemes, namely the BMSN scheme as proposed in , and in the CSN scheme as proposed in . For that purpose we have constructed in a five-flavor parton density set from a four-flavor parton density set. Using our densities we will study the differences between the bottom components of the deep inelastic structure functions $`F_{i,b}^{\mathrm{CSN}}(n_f+1)`$ and $`F_{i,b}^{\mathrm{BMSN}}(n_f+1)`$, where the number of light flavors is $`n_f=4`$.
To keep the discussion short we simply refer the reader to Sec.III of for a discussion of the $`\overline{\mathrm{MS}}`$ parton densities, the exact solution for the running coupling constant and the scale choice. All references to three-flavour (four-flavor) densities should be replaced by four-flavor (five-flavor) densities respectively. All our calculations of next-to-leading (NLO) and next-to-next-leading order (NNLO) quantities use $`\mathrm{\Lambda }_{3,4,5,6}^{\overline{\mathrm{MS}}}=299.4,246,167.7,67.8\mathrm{MeV}`$, which yields $`\alpha _s(5,M_Z^2)=0.114`$. The structure functions are defined in Eqs.(3.9)-(3.17) of , where now $`n_f=4`$ and $`m_c`$ is replaced by $`m_b=4.5`$ $`\mathrm{GeV}/\mathrm{c}^2`$. To make this paper reasonably self-contained we now reproduce the final formulae we use for the structure functions. For $`i=2,L`$ the CSN scheme uses
$`F_{i,Q}^{\mathrm{CSN}}(n_f+1,\mathrm{\Delta },Q^2,m^2)=e_Q^2[f_{Q+\overline{Q}}^{\mathrm{NNLO}}(n_f+1,\mu ^2)𝒞_{i,Q}^{\mathrm{CSN},\mathrm{NS},(0)}({\displaystyle \frac{Q^2}{m^2}})`$ (II.1)
$`+a_s(n_f+1,\mu ^2)\{f_{Q+\overline{Q}}^{\mathrm{NLO}}(n_f+1,\mu ^2)𝒞_{i,Q}^{\mathrm{CSN},\mathrm{NS},(1)}({\displaystyle \frac{Q^2}{m^2}},{\displaystyle \frac{Q^2}{\mu ^2}})`$ (II.2)
$`+f_g^{\mathrm{S},\mathrm{NLO}}(n_f+1,\mu ^2)𝒞_{i,g}^{\mathrm{CSN},\mathrm{S},(1)}({\displaystyle \frac{Q^2}{m^2}},{\displaystyle \frac{Q^2}{\mu ^2}})\}`$ (II.3)
$`+a_s^2(n_f+1,\mu ^2)\{f_{Q+\overline{Q}}^{\mathrm{LO}}(n_f+1,\mu ^2)(𝒞_{i,q}^{\mathrm{NS},(2)}(n_f+1,{\displaystyle \frac{Q^2}{m^2}},{\displaystyle \frac{Q^2}{\mu ^2}})`$ (II.4)
$`+𝒞_{i,q}^{\mathrm{PS},(2)}({\displaystyle \frac{Q^2}{m^2}},{\displaystyle \frac{Q^2}{\mu ^2}}))+{\displaystyle }_{l=1}^{n_f}f_{l+\overline{l}}^{\mathrm{LO}}(n_f+1,\mu ^2)𝒞_{i,q}^{\mathrm{CSN},\mathrm{PS},(2)}({\displaystyle \frac{Q^2}{m^2}},{\displaystyle \frac{Q^2}{\mu ^2}})`$ (II.5)
$`+f_g^{\mathrm{S},\mathrm{LO}}(n_f+1,\mu ^2)𝒞_{i,g}^{\mathrm{CSN},\mathrm{S},(2)}({\displaystyle \frac{Q^2}{m^2}},{\displaystyle \frac{Q^2}{\mu ^2}})\}]`$ (II.6)
$`+a_s^2(n_f+1,\mu ^2){\displaystyle \underset{k=1}{\overset{n_f}{}}}e_k^2f_{k+\overline{k}}^{\mathrm{LO}}(n_f+1,\mu ^2)L_{i,q}^{\mathrm{HARD},\mathrm{NS},(2)}(\mathrm{\Delta },{\displaystyle \frac{Q^2}{m^2}}).`$ (II.7)
where we choose the heavy quark $`Q`$ to be the bottom quark, so that the number of light flavors is $`n_f=4`$. The charge of the bottom quark is $`e_Q=1/3`$ and its mass is $`m=m_b`$. The coefficient function $`L_{i,q}^{\mathrm{HARD}}`$ depends on a parameter $`\mathrm{\Delta }`$ which refers to the invariant mass of the $`Q\overline{Q}`$-pair. For bottom quark production we choose $`\mathrm{\Delta }=100`$ $`(\mathrm{GeV}/c)^2`$. The coefficient functions labelled by $`C^{\mathrm{CSN}}`$ depend on the heavy quark mass but are finite in the limit $`m0`$. They are defined in Eqs.(2.8)-(2.20) in . To simplify the notation we will refer to the above structure functions as $`F_{i,b}^{\mathrm{CSN}}(n_f=5)`$, indicating that they depend on five-flavor parton densities.
The same parameters $`e_Q`$, $`m_b`$, $`\mathrm{\Delta }`$ etc., also show up in the expressions for the bottom quark structure functions in the BMSN scheme. Here we have the representations
$`F_{i,Q}^{\mathrm{BMSN}}(n_f+1,\mathrm{\Delta },Q^2,m^2)=F_{i,Q}^{\mathrm{EXACT}}(n_f,\mathrm{\Delta },Q^2,m^2)`$ (II.8)
$`F_{i,Q}^{\mathrm{ASYMP}}(n_f,\mathrm{\Delta },Q^2,m^2)+F_{i,Q}^{\mathrm{PDF}}(n_f+1,\mathrm{\Delta },Q^2,m^2).`$ (II.9)
The pieces in this formulae represent first the results in FOPT, given by
$`F_{i,Q}^{\mathrm{EXACT}}(n_f,\mathrm{\Delta },Q^2,m^2)=e_Q^2[a_s(n_f,\mu ^2)f_g^{\mathrm{S},\mathrm{NLO}}(n_f,\mu ^2)H_{i,g}^{\mathrm{S},(1)}\left({\displaystyle \frac{Q^2}{m^2}}\right)`$ (II.10)
$`+a_s^2(n_f,\mu ^2)\{{\displaystyle \underset{k=1}{\overset{n_f}{}}}f_{k+\overline{k}}^{\mathrm{LO}}(n_f,\mu ^2)H_{i,q}^{\mathrm{PS},(2)}({\displaystyle \frac{Q^2}{m^2}},{\displaystyle \frac{Q^2}{\mu ^2}})`$ (II.11)
$`+f_g^{\mathrm{S},\mathrm{LO}}(n_f,\mu ^2)H_{i,g}^{\mathrm{S},(2)}({\displaystyle \frac{Q^2}{m^2}},{\displaystyle \frac{Q^2}{\mu ^2}})\}]`$ (II.12)
$`+a_s^2(n_f,\mu ^2){\displaystyle \underset{k=1}{\overset{n_f}{}}}e_k^2f_{k+\overline{k}}^{\mathrm{LO}}(n_f,\mu ^2)L_{i,q}^{\mathrm{HARD},\mathrm{NS},(2)}(\mathrm{\Delta },{\displaystyle \frac{Q^2}{m^2}}).`$ (II.13)
The next pieces are the structure functions $`F_{i,Q}^{\mathrm{ASYMP}}(n_f)`$ which can be obtained from $`F_{i,Q}^{\mathrm{EXACT}}(n_f)`$ by replacing all exact heavy quark coefficient functions $`H_{i,k}`$ and $`L_{i,k}`$ ($`k=q,g`$) by their asymptotic analogues which are defined by
$`H_{i,k}^{\mathrm{ASYMP}}=\underset{Q^2m^2}{\text{lim}}H_{i,k},L_{i,k}^{\mathrm{ASYMP}}=\underset{Q^2m^2}{\text{lim}}L_{i,k}.`$ (II.14)
Finally the structure functions $`F_{i,Q}^{\mathrm{PDF}}(n_f+1)`$ which are very often called the ZM-VFNS representations are defined by
$`F_{i,Q}^{\mathrm{PDF}}(n_f+1,\mathrm{\Delta },Q^2,m^2)=e_Q^2[f_{Q+\overline{Q}}^{\mathrm{NNLO}}(n_f+1,\mu ^2)𝒞_{i,q}^{\mathrm{NS},(0)}`$ (II.15)
$`+a_s(n_f+1,\mu ^2)\{f_{Q+\overline{Q}}^{\mathrm{NLO}}(n_f+1,\mu ^2)𝒞_{i,q}^{\mathrm{NS},(1)}({\displaystyle \frac{Q^2}{\mu ^2}})`$ (II.16)
$`+f_g^{\mathrm{S},\mathrm{NLO}}(n_f+1,\mu ^2)\stackrel{~}{𝒞}_{i,g}^{\mathrm{S},(1)}({\displaystyle \frac{Q^2}{\mu ^2}})\}`$ (II.17)
$`+a_s^2(n_f+1,\mu ^2)\{f_{Q+\overline{Q}}^{\mathrm{LO}}(n_f+1,\mu ^2)(𝒞_{i,q}^{\mathrm{NS},(2)}(n_f+1,{\displaystyle \frac{Q^2}{\mu ^2}})`$ (II.18)
$`+\stackrel{~}{𝒞}_{i,q}^{\mathrm{PS},(2)}\left({\displaystyle \frac{Q^2}{\mu ^2}}\right))+{\displaystyle \underset{l=1}{\overset{n_f}{}}}f_{l+\overline{l}}^{\mathrm{LO}}(n_f+1,\mu ^2))\stackrel{~}{𝒞}_{i,q}^{\mathrm{PS},(2)}({\displaystyle \frac{Q^2}{\mu ^2}})`$ (II.19)
$`+f_g^{\mathrm{S},\mathrm{LO}}(n_f+1,\mu ^2)\stackrel{~}{𝒞}_{i,g}^{\mathrm{S},(2)}({\displaystyle \frac{Q^2}{\mu ^2}})\}]`$ (II.20)
$`+a_s^2(n_f+1,\mu ^2){\displaystyle \underset{k=1}{\overset{n_f}{}}}e_k^2f_{k+\overline{k}}^{\mathrm{LO}}(n_f,\mu ^2)L_{i,q}^{\mathrm{HARD},\mathrm{ASYMP},\mathrm{NS},(2)}(\mathrm{\Delta },{\displaystyle \frac{Q^2}{m^2}}).`$ (II.21)
In all these results the heavy quark Q refers to the bottom quark and the other parameters are defined above. For simplicity we will refer to these structure functions as $`F_{i,b}^{\mathrm{EXACT}}(n_f=4)`$, $`F_{i,b}^{\mathrm{ASYMP}}(n_f=4)`$, and $`F_{i,b}^{\mathrm{PDF}}(n_f=5)`$ respectively, which indicates that the first two structure functions depend on four-flavor densities and the last one depends on five-flavor densities. The parton densities $`f_k`$ in the above formulae are represented in leading order (LO), next-to-leading order (NLO) and next-to-next-to-leading order (NNLO). The NNLO case refers to the boundary conditions imposed in since the three-loop splitting functions are not known yet. These parton densities have been constructed in starting from a three-flavor parametrization in . The multiplication of the densities with the heavy and light parton coefficient functions is done in such a way that the perturbation series is strictly truncated at order $`\alpha _s^2`$. This is necessary to avoid scheme dependent terms which would otherwise arise beyond order $`\alpha _s^2`$. Therefore the following requirement is satisfied
$`F_{i,Q}^{\mathrm{CSN}}(n_f=5)=F_{i,Q}^{\mathrm{BMSN}}(n_f=5)=F_{i,Q}^{\mathrm{EXACT}}(n_f=4)\text{for}Q^2m^2.`$ (II.23)
Since $`f_Q(m^2)^{\mathrm{NNLO}}0`$ (see ) this condition can be only satisfied when we truncate the perturbation series at the same order. Furthermore because of Eq. (II.14) and the property
$`\underset{Q^2m^2}{\text{lim}}𝒞_{i,k}^{\mathrm{CSN},(\mathrm{l})}(n_f+1,{\displaystyle \frac{Q^2}{m^2}},{\displaystyle \frac{Q^2}{\mu ^2}})=𝒞_{i,k}^{(l)}(n_f+1,{\displaystyle \frac{Q^2}{\mu ^2}}),`$ (II.24)
we have the asymptotic relation
$`\underset{Q^2m^2}{\text{lim}}F_{i,Q}^{\mathrm{BMSN}}(n_f+1,\mathrm{\Delta },Q^2,m^2)=\underset{Q^2m^2}{\text{lim}}F_{i,Q}^{\mathrm{CSN}}(n_f+1,\mathrm{\Delta },Q^2,m^2)`$ (II.25)
$`=F_{i,Q}^{\mathrm{PDF}}(n_f+1,\mathrm{\Delta },Q^2,m^2).`$ (II.26)
At first sight the form of the expression for $`F_{i,Q}^{\mathrm{BMSN}}`$ in Eq. (II.8) looks quite different from the one presented for $`F_{i,Q}^{\mathrm{CSN}}`$ in Eq. (II.1). However this is not true. Using the mass factorization relations for the asymptotic heavy quark coefficient functions in one can cast $`F_{i,Q}^{\mathrm{BMSN}}`$ into the same form as presented for $`F_{i,Q}^{\mathrm{CSN}}`$ where all quark coefficient functions of the type $`𝒞_{i,Q}^{\mathrm{CSN}}`$ are replaced by their light quark analogues $`𝒞_{i,q}`$ appearing in Eq. (II.15). This replacement also applies to the $`𝒞_{i,Q}^{\mathrm{CSN}}`$ occurring in the mass factorization relations for $`𝒞_{i,g}^{\mathrm{CSN},\mathrm{S}}`$ and $`𝒞_{i,q}^{\mathrm{CSN},\mathrm{PS}}`$ presented in . Therefore the difference between the CSN and BMSN schemes can be attributed to the powers $`(m^2/Q^2)^j`$ showing up in $`𝒞_{i,Q}^{\mathrm{CSN}}`$ but absent in $`𝒞_{i,q}`$. This effect is only noticeable in the threshold region where $`Q^2m^2`$ as we will show below.
The heavy quark coefficient functions $`𝒞_{i,k}^{\mathrm{CSN}}`$, $`H_{i,k}`$, $`L_{i,k}`$ ($`k=Q,q,g`$) and the light partonic coefficient functions $`𝒞_{i,k}`$ ($`k=q,g`$) can be traced back to the following processes
$`𝒞_{i,g}^{\mathrm{CSN},\mathrm{S},(1)},H_{i,g}^{\mathrm{S},(1)}`$ $`:`$ $`\gamma ^{}+gQ+\overline{Q}\text{[6]}\text{ (CSN), }\text{[1]}\text{ (EXACT),}`$ (II.28)
(ASYMP)
$`𝒞_{i,g}^{\mathrm{CSN},\mathrm{S},(2)},H_{i,g}^{\mathrm{S},(2)}`$ $`:`$ $`\gamma ^{}+gQ+\overline{Q}+g\text{[6]}\text{ (CSN), }\text{[1]}\text{ (EXACT),}`$ (II.30)
(ASYMP)
$`𝒞_{i,q}^{\mathrm{CSN},\mathrm{PS},(2)},H_{i,q}^{\mathrm{PS},(2)}`$ $`:`$ $`\gamma ^{}+qQ+\overline{Q}+q\text{Bethe-Heitler reaction}`$ (II.32)
(CSN), (EXACT), (ASYMP)
$`L_{i,q}^{\mathrm{HARD},\mathrm{NS},(2)}`$ $`:`$ $`\gamma ^{}+qQ+\overline{Q}+q\text{Compton reaction}`$ (II.34)
(EXACT and ASYMP)
$`𝒞_{i,Q}^{\mathrm{CSN},\mathrm{NS},(0)},H_{i,Q}^{\mathrm{NS},(0)}`$ $`:`$ $`\gamma ^{}+QQ`$ (II.35)
$`𝒞_{i,Q}^{\mathrm{CSN},\mathrm{NS},(1)},H_{i,Q}^{\mathrm{NS},(1)}`$ $`:`$ $`\gamma ^{}+QQ+g\text{[12]}`$ (II.36)
$`𝒞_{i,q}^{\mathrm{NS},(0)}`$ $`:`$ $`\gamma ^{}+qq`$ (II.37)
$`𝒞_{i,q}^{\mathrm{NS},(1)}`$ $`:`$ $`\gamma ^{}+qq+g\text{[13]}`$ (II.38)
$`𝒞_{i,q}^{\mathrm{NS},(2)}`$ $`:`$ $`\gamma ^{}+qq+g+g\text{[13]}`$ (II.39)
$`𝒞_{i,q}^{\mathrm{NS},(2)}`$ $`:`$ $`\gamma ^{}+qq+\overline{q}+q\text{[13]}`$ (II.40)
$`\stackrel{~}{𝒞}_{i,q}^{\mathrm{PS},(2)}`$ $`:`$ $`\gamma ^{}+qq+\overline{q}+q\text{[13]}`$ (II.41)
$`\stackrel{~}{𝒞}_{i,g}^{\mathrm{S},(1)}`$ $`:`$ $`\gamma ^{}+gq+\overline{q}\text{[13]}`$ (II.42)
$`\stackrel{~}{𝒞}_{i,g}^{\mathrm{S},(2)}`$ $`:`$ $`\gamma ^{}+gq+\overline{q}+g\text{[13]}.`$ (II.43)
Behind the reactions we have quoted the references in which the corresponding coefficient functions can be found. Note that the heavy quark coefficient functions $`H_{i,k}`$ are mass singular when $`m0`$. This can be observed immediately when one looks at $`H_{i,k}^{\mathrm{ASYMP}}`$ which behaves like $`\mathrm{ln}^m(\mu ^2/m^2)\mathrm{ln}^n(Q^2/m^2)`$ (see ). After the logarithms are removed one obtains the quantities $`𝒞_{i,k}^{\mathrm{CNS}}`$ which, even though they depend on $`m`$, are finite in the limit $`m0`$. The coefficient function $`L_{i,k}^{\mathrm{HARD}}`$ is finite by itself because as we mentioned above we have imposed a lower cut off $`\mathrm{\Delta }=100`$ $`(\mathrm{GeV}/c)^2`$ on the invariant mass of the $`Q\overline{Q}`$-pair. Finally notice that all parton densities, coefficient functions and the running coupling constant are presented in the $`\overline{\mathrm{MS}}`$-scheme.
Now we present results for the various structure functions. We are interested in the bottom quark structure functions $`F_{i,b}^{\mathrm{CSN}}(n_f=5)`$ and $`F_{i,b}^{\mathrm{BMSN}}(n_f=5)`$ for $`i=2,L`$ in NNLO for the CSN and BMSN schemes respectively. In Fig. 1 we have plotted the structure functions $`F_{2,b}^{\mathrm{CSN}}(n_f=5)`$, $`F_{2,b}^{\mathrm{BMSN}}(n_f=5)`$, $`F_{2,b}^{\mathrm{PDF}}(n_f=5)`$ and $`F_{2,b}^{\mathrm{EXACT}}(n_f=4)`$ in the region $`20<Q^2<10^3`$ in units of $`(\mathrm{GeV}/\mathrm{c})^2`$ for $`x=0.05`$. This figure reveals that there is hardly any difference between the BMSN and CSN prescriptions. The curves in both prescriptions are essentially identical to that for $`F_{2,b}^{\mathrm{EXACT}}(n_f=4)`$. In this region $`F_{2,b}^{\mathrm{PDF}}(n_f=5)`$ is larger than the other results which is expected from the discussion of the bottom quark density given in . There is still an appreciable difference at the highest plotted $`Q^2`$ demonstrating that mass effects are important up to very large scales. Notice that for $`Q^235`$ $`(\mathrm{GeV}/\mathrm{c})^2`$ $`F_{2,b}^{\mathrm{PDF}}(n_f=5)`$ becomes negative which means that bottom quark electroproduction cannot be described by this quantity anymore. In Fig. 2 we present the same plots for $`x=0.005`$. Again one cannot distinguish between $`F_{2,b}^{\mathrm{BMSN}}(n_f=5)`$ and $`F_{2,c}^{\mathrm{CSN}}(n_f=5)`$ but now both are smaller than $`F_{2,b}^{\mathrm{EXACT}}(n_f=4)`$ over the whole $`Q^2`$ range. The latter is smaller than $`F_{2,b}^{\mathrm{PDF}}(n_f=5)`$ in particular for $`Q^2>35(\mathrm{GeV}/\mathrm{c})^2`$. Further we want to emphasize that due to our careful treatment of the threshold region there is an excellent cancellation between $`F_{2,b}^{\mathrm{PDF}}(n_f=5)`$ and $`F_{2,b}^{\mathrm{ASYMP}}(n_f=4)`$ so that both $`F_{2,b}^{\mathrm{CSN}}(n_f=5)`$ and $`F_{2,b}^{\mathrm{BMSN}}(n_f=5)`$ tend to $`F_{2,b}^{\mathrm{EXACT}}(n_f=4)`$ at $`Q^2=m_b^2`$. At large $`Q^2`$ we have a cancellation between $`F_{2,b}^{\mathrm{ASYMP}}(n_f=4)`$ and $`F_{2,b}^{\mathrm{EXACT}}(n_f=4)`$ so that both $`F_{2,b}^{\mathrm{CSN}}(n_f=5)`$ and $`F_{2,b}^{\mathrm{BMSN}}(n_f=5)`$ slowly tend to $`F_{2,b}^{\mathrm{PDF}}(n_f=5)`$. They are only identical at extremely large $`Q^2`$ demonstrating that mass effects are still important over a wide range in $`x`$ and $`Q^2`$.
In Fig.3 we show similar plots as in Fig.1 for the bottom quark longitudinal structure functions. Here we observe a small difference between the plots for $`F_{L,b}^{\mathrm{CSN}}(n_f=5)`$ and $`F_{L,b}^{\mathrm{BMSN}}(n_f=5)`$ in the region $`20<Q^2<10^3(\mathrm{GeV}/\mathrm{c})^2`$. Furthermore $`F_{L,b}^{\mathrm{PDF}}(n_f=5)`$ is considerably larger than the other three structure functions, which differs from the behavior seen in Fig.1. This can be mainly attributed to the gluon density which plays a more prominant role in $`F_{L,b}`$ than in $`F_{2,b}`$. For $`x=0.005`$ (see Fig.4) the small difference between the BMSN and the CSN descriptions becomes more conspicuous for low $`Q^2`$.
In Figs.5 and 6 we make a comparison between the NLO and the NNLO structure functions $`F_{2,b}^{\mathrm{CSN}}(n_f=5)`$ and $`F_{2,b}^{\mathrm{BMSN}}(n_f=5)`$. Both the CSN and and BMSN descriptions lead to the same results in both NLO and NNLO. However while going from NLO to NNLO the the structure functions $`F_{2,b}^{\mathrm{CSN}}(n_f=5)`$ and $`F_{2,b}^{\mathrm{BMSN}}(n_f=5)`$ increase a little bit. The differences in the case of $`x=0.005`$ in Fig.6 are even smaller than those observed for $`x=0.05`$ in Fig.5. The same comparison between NLO and NNLO results is made for the longitudinal structure functions in Figs.7 and 8. Here the differences between NLO and NNLO cases are much larger than in the case of $`F_{2,b}`$ in Figs.5,6. In NLO both $`F_{L,b}^{\mathrm{BMSN}}(n_f=5)`$ and $`F_{L,b}^{\mathrm{CSN}}(n_f=5)`$ are smaller than the NNLO results.
Previous results for $`F_{i,b}^{\mathrm{EXACT}}(x,Q^2,m_b^2)`$ and have been presented in Figs. 20a,20b in for a now obsolete set of parton densities, so the values quoted there are too small. To show these changes we add in Figs.9 and 10 plots for the $`x`$ dependence of $`F_{2,b}^{\mathrm{CSN}}(n_f=5)`$, $`F_{2,b}^{\mathrm{BMSN}}(n_f=5)`$, $`F_{2,b}^{\mathrm{PDF}}(n_f=5)`$ and $`F_{2,b}^{\mathrm{EXACT}}(n_f=4)`$ at fixed $`Q^2=30`$ and $`Q^2=100`$ in units of $`(\mathrm{GeV}/\mathrm{c})^2`$ respectively. Finally we also show in Figs.11 and 12 plots for the $`x`$ dependence of $`F_{L,b}^{\mathrm{CSN}}(n_f=5)`$, $`F_{L,b}^{\mathrm{BMSN}}(n_f=5)`$, $`F_{L,b}^{\mathrm{PDF}}(n_f=5)`$ and $`F_{L,b}^{\mathrm{EXACT}}(n_f=4)`$ at fixed $`Q^2=30`$ and $`Q^2=100`$ in units of $`(\mathrm{GeV}/\mathrm{c})^2`$ respectively. Note that there are also some recent results for $`F_{2,b}`$ in in the TR scheme and in for FOPT.
The plots for $`F_{L,b}^{\mathrm{CSN}}(n_f=5)`$ in Figs.7,8 do not show a negative region at small $`Q^2`$ which could have been expected by analogy with the results for $`F_{L,c}^{\mathrm{CSN}}(n_f=4)`$ in . In the case of bottom quark production the negative regions do occur but at even smaller values of $`x`$. In Figs. 13, 14 we show the same plots as in Figs.1,3 respectively but for $`x=5\times 10^5`$. Now the structure function $`F_{L,b}^{\mathrm{CSN}}`$ is negative in the region $`Q^230`$ $`(\mathrm{GeV}/\mathrm{c})^2`$. In Figs. 15,16 we show the same plots as in Figs. 5,7 respectively but for $`x=5\times 10^5`$. Fig. 16 shows that the longitudinal structure functions for the case of bottom production also have negative regions at small $`Q^2`$ values in both NLO and NNLO. This phenomenon also occurs for the charm quark structure functions in . In the NLO case this arises because the term $`f_g^{\mathrm{S},\mathrm{NLO}}(n_f+1,\mu ^2)𝒞_{L,g}^{\mathrm{CSN},\mathrm{S},(1)}(Q^2/m^2,Q^2/\mu ^2)`$ in Eq. (II.1) is negative due to the definition of the gluon coefficient function in the CSN scheme (see Eq.(2.19)) in ) which is given by
$`𝒞_{L,g}^{\mathrm{CSN},\mathrm{S},(1)}({\displaystyle \frac{Q^2}{m^2}},{\displaystyle \frac{Q^2}{\mu ^2}})`$ $`=`$ $`H_{L,g}^{\mathrm{S},(1)}\left({\displaystyle \frac{Q^2}{m^2}}\right)A_{Qg}^{\mathrm{S},(1)}({\displaystyle \frac{\mu ^2}{m^2}})𝒞_{L,Q}^{\mathrm{CSN},\mathrm{NS},(0)}\left({\displaystyle \frac{Q^2}{m^2}}\right),`$ (II.44)
$`\text{with}𝒞_{L,Q}^{\mathrm{CSN},\mathrm{NS},(0)}`$ $`=`$ $`{\displaystyle \frac{4m^2}{Q^2}},`$ (II.45)
where $`A_{Qg}^{\mathrm{S},(1)}`$ denotes the one-loop operator matrix element computed in . Notice that the latter and the lowest order exact coefficient function $`H_{L,g}^{\mathrm{S},(1)}`$ are always positive. Because of the minus sign in Eq. (II.44) it appears that the coefficient function $`𝒞_{L,g}^{\mathrm{CSN},\mathrm{S},(1)}`$ can become negative in particular at low $`Q^2`$ values. In the NNLO case one obtains more negative contributions due to the term $`f_{Q+\overline{Q}}^{\mathrm{NNLO}}(n_f+1,\mu ^2)𝒞_{L,Q}^{\mathrm{CSN},\mathrm{NS},(0)}`$ in formula (II.1). It turns out that $`f_{Q+\overline{Q}}^{\mathrm{NNLO}}(n_f+1,x,\mu ^2)`$ is negative at small $`x`$ and $`\mu ^2=Q^2m^2`$. Notice that at the latter scale $`f_{Q+\overline{Q}}^{\mathrm{LO}}(n_f+1,x,\mu ^2)`$ and $`f_{Q+\overline{Q}}^{\mathrm{NLO}}(n_f+1,x,\mu ^2)`$ are very small because they vanish at $`\mu =m`$ in contrast to $`f_{Q+\overline{Q}}^{\mathrm{NNLO}}(n_f+1,x,\mu ^2)`$. The behavior of the structure function above is characteristic of the CSN scheme since it does not appear in the case of BMSN. This is because in the latter scheme the longitudinal coefficient function, represented by $`𝒞_{L,q}^{\mathrm{CSN},\mathrm{NS},(0)}`$, is identical to zero so that the zeroth order contribution to $`F_{L,b}^{\mathrm{BMSN}}(n_f=5)`$ vanishes and the first order correction is given by $`𝒞_{L,g}^{\mathrm{BMSN},\mathrm{S},(1)}=H_{L,g}^{\mathrm{S},(1)}`$. The latter leads to a positive structure function over the whole kinematical region. To further demonstrate this point we plot in Fig. 17 pieces of the NLO result
$`F_{L,b}^{\mathrm{CSN}}(n_f+1,Q^2,m^2)=e_b^2[f_{b+\overline{b}}^{\mathrm{NLO}}(n_f+1,\mu ^2)𝒞_{L,b}^{\mathrm{CSN},\mathrm{NS},(0)}({\displaystyle \frac{Q^2}{m^2}})`$ (II.46)
$`+a_s(n_f+1,\mu ^2)\{f_{b+\overline{b}}^{\mathrm{LO}}(n_f+1,\mu ^2)𝒞_{L,b}^{\mathrm{CSN},\mathrm{NS},(1)}({\displaystyle \frac{Q^2}{m^2}},{\displaystyle \frac{Q^2}{\mu ^2}})`$ (II.47)
$`+f_g^{\mathrm{S},\mathrm{LO}}(n_f+1,\mu ^2)𝒞_{L,g}^{\mathrm{CSN},\mathrm{S},(1)}({\displaystyle \frac{Q^2}{m^2}},{\displaystyle \frac{Q^2}{\mu ^2}})\}.`$ (II.48)
The sum of the b-quark contributions, labelled Term1, is always positive. The gluonic contribution, labelled Term2, is clearly negative over a wide range in $`Q^2`$ and large enough that the sum $`F_{L,b}^{\mathrm{CSN}}`$ is also negative for $`30Q^2150`$ $`\mathrm{GeV}^2`$, as in Fig.16. This behavior is due to the gluonic coefficient function $`𝒞_{L,g}^{\mathrm{CSN},\mathrm{S},(1)}`$ which is explained below Eq. (II.44). However the magnitude of the gluonic contribution depends on the choice of the gluon density. If we use an NLO gluon density in the order $`\alpha _s`$ contribution to the structure function in Eq. (II.46) rather than a LO gluon density then the sum of the first two terms is unchanged but the gluonic part is now smaller in magnitude. These contributions are shown in Fig.18 where now $`|\mathrm{Term2}|\mathrm{Term1}`$ so that the total result for $`F_{L,b}^{\mathrm{CSN}}`$ is everywhere positive. However this procedure violates our prescription for the computation of the structure functions in both the CSN and the BMSN schemes. In this prescription the LO densities are multiplied by the highest order coefficient function whereas the NLO densities are combined with lower order coefficient functions (see formulae (II.1), (II.10) and (II.15)). In this way the perturbation series is truncated up to the order we want to compute the structure functions. Hence we avoid terms, arising beyond that order, which introduce a scheme dependence and spoils the threshold behavior (see ). The latter happens if one follows the usual procedure where one multiplies the highest order densities by the highest order coefficient functions. The difference between the usual procedure and our prescription is not only shown by our parton density set but is also observed for other sets presented in the literature. Examples are recent sets such as MRST98 (with $`m_b=4.3`$ $`\mathrm{GeV}`$, $`m_c=1.35`$ $`\mathrm{GeV}`$), MRST99 (with $`m_b=4.3`$ $`\mathrm{GeV}`$, $`m_c=1.43`$ $`\mathrm{GeV}`$), and CTEQ5 (with $`m_b=4.5`$ $`\mathrm{GeV}`$, $`m_c=1.3`$ $`\mathrm{GeV}`$). Note that the MRST99 set does not provide LO densities. Using their NLO densities they yield positive values for the $`Q^2`$ dependence of $`F_{L,b}^{\mathrm{CSN}}`$ at $`x=5\times 10^5`$. There are both LO (CTEQ5L) and NLO (CTEQ5M) densities in the CTEQ5 set and we have checked that, for the same $`x`$, $`Q^2`$ values $`F_{L,b}^{\mathrm{CSN}}`$ is positive with purely NLO densities but has a negative region when the LO and NLO densities are used according to our prescription. The observations made above leads to the conclusion that the $`4m^2/Q^2`$ term in the non-singlet CSN longitudinal coefficient function in Eq. (II.44) leads to a negative gluonic coefficient function. When the latter is used together with the latest LO and NLO parton density sets $`F_{L,b}^{\mathrm{CSN}}(n_f=5)`$ becomes negative in the low $`Q^2`$-region at small $`x`$. We speculate that the CSN scheme would always yield positive structure functions if we could use parton density sets which fitted data with convolutions of LO densities with $`O(\alpha _s)`$ coefficient functions and NLO densities with zeroth order coefficient functions. Unfortunately such densities are not available.
To summarize the main points we have implemented two variable flavor number schemes for bottom quark electroproduction in NNLO and compared them with NLO FOPT results. The schemes differ in the way mass factorization is implemented. In the CSN scheme this is done with respect to the full heavy and light quark structure functions at finite $`Q^2`$. In the BMSN scheme the mass factorization is only applied to the coefficient functions in the large $`Q^2`$ limit. Both schemes require four-flavor and five-flavor parton densities which satisfy discontinuous NNLO matching conditions at a scale $`\mu =m_b`$. We have constructed these densities using our own evolution code . The schemes also require matching conditions on the coefficient functions which are implemented in this paper. Note that we have also removed the dangerous terms in $`\mathrm{ln}^3(Q^2/m_b^2)`$ from the Compton contributions so that both $`F_{i,b}^{\mathrm{CSN}}(n_f=5)`$ and $`F_{i,b}^{\mathrm{BMSN}}(n_f=5)`$ are collinear safe. As in we have done this in a way which is consistent with our study of inclusive quantities by implementing a cut $`\mathrm{\Delta }`$ on the mass of the $`b\overline{b}`$ pair. We stress that any ZM-VFNS bottom quark density description of $`F_{i,b}`$ must use collinear safe definitions. This is not required in the fixed order perturbation theory approach given by $`F_{i,b}^{\mathrm{EXACT}}(n_f=4)`$ in for moderate $`Q^2`$-values.
Finally we made a careful analysis of the threshold behaviors of $`F_{i,b}^{\mathrm{CSN}}(n_f=5)`$ and $`F_{i,b}^{\mathrm{BMSN}}(n_f=5)`$. In order to achieve the required cancellations at the scale $`\mu =m_b`$ so that they both become equal to $`F_{i,b}^{\mathrm{EXACT}}(n_f=4)`$ one must be very careful to combine terms with the same order in the expansion in $`\alpha _s`$. The approximation we made in this paper, of using NLO splitting functions in place of NNLO splitting functions, was sufficient for our purposes. We successfully implemented the required cancellations near threshold and the corresponding limits at large scales came out naturally. Inconsistent sets of parton densities automatically spoil these cancelations. Since there are only minor differences between the CSN, BMSN and NLO FOPT predictions it is clear that the use of variable flavor number schemes for bottom quark production is not required for the analysis of HERA data. However the ZM-VFNS description is clearly inadequate at small $`Q^2`$.
ACKNOWLEDGMENTS
The work of A. Chuvakin and J. Smith was partially supported by the National Science Foundation grant PHY-9722101. The work of W.L. van Neerven was supported by the EC network ‘QCD and Particle Structure’ under contract No. FMRX–CT98–0194.
Figure Captions
The bottom quark structure functions $`F_{2,b}^{\mathrm{EXACT}}(n_f=4)`$ (solid line) $`F_{2,b}^{\mathrm{CSN}}(n_f=5)`$, (dot-dashed line) $`F_{2,b}^{\mathrm{BMSN}}(n_f=5)`$, (dashed line) and $`F_{2,b}^{\mathrm{PDF}}(n_f=5)`$, (dotted line) in NNLO for $`x=0.05`$ plotted as functions of $`Q^2`$.
Same as in Fig. 1 but now for $`x=0.005`$.
The bottom quark structure functions $`F_{L,b}^{\mathrm{EXACT}}(n_f=4)`$ (solid line) $`F_{L,b}^{\mathrm{CSN}}(n_f=5)`$, (dot-dashed line) $`F_{L,b}^{\mathrm{BMSN}}(n_f=5)`$, (dashed line) and $`F_{L,b}^{\mathrm{PDF}}(n_f=5)`$, (dotted line) in NNLO for $`x=0.05`$ plotted as functions of $`Q^2`$.
Same as in Fig. 3 but now for $`x=0.005`$.
The bottom quark structure functions $`F_{2,b}^{\mathrm{BMSN}}(n_f=5)`$ in NLO (solid line), NNLO (dotted line) and $`F_{2,b}^{\mathrm{CSN}}(n_f=5)`$ in NLO (dashed line), NNLO (dot-dashed line) for $`x=0.05`$ plotted as functions of $`Q^2`$.
Same as in Fig. 5 but now for $`x=0.005`$.
The bottom quark structure functions $`F_{L,b}^{\mathrm{BMSN}}(n_f=5)`$ in NLO (solid line), NNLO (dotted line) and $`F_{L,b}^{\mathrm{CSN}}(n_f=5)`$ in NLO (dashed line), NNLO (dot-dashed line) for $`x=0.05`$ plotted as functions of $`Q^2`$.
Same as in Fig. 7 but now for $`x=0.005`$.
The bottom quark structure functions $`F_{2,b}^{\mathrm{EXACT}}(n_f=4)`$ (solid line) $`F_{2,b}^{\mathrm{CSN}}(n_f=5)`$, (dot-dashed line) $`F_{2,b}^{\mathrm{BMSN}}(n_f=5)`$, (dashed line) and $`F_{2,b}^{\mathrm{PDF}}(n_f=5)`$, (dotted line) in NNLO for $`Q^2=30`$ $`(\mathrm{GeV}/\mathrm{c})^2`$ plotted as functions of $`x`$.
Same as in Fig. 9 but now for $`Q^2=100`$ $`(\mathrm{GeV}/\mathrm{c})^2`$.
The bottom quark structure functions $`F_{L,b}^{\mathrm{EXACT}}(n_f=4)`$ (solid line) $`F_{L,b}^{\mathrm{CSN}}(n_f=5)`$, (dot-dashed line) $`F_{L,b}^{\mathrm{BMSN}}(n_f=5)`$, (dashed line) and $`F_{L,b}^{\mathrm{PDF}}(n_f=5)`$, (dotted line) in NNLO for $`Q^2=30`$ $`(\mathrm{GeV}/\mathrm{c})^2`$ plotted as functions of $`x`$.
Same as in Fig. 11 but now for $`Q^2=100`$ $`(\mathrm{GeV}/\mathrm{c})^2`$.
Same as in Fig. 1 but now for $`x=5\times 10^5`$.
Same as in Fig. 3 but now for $`x=5\times 10^5`$.
Same as in Fig. 5 but now for $`x=5\times 10^5`$.
Same as in Fig. 7 but now for $`x=5\times 10^5`$.
The bottom quark structure function $`F_{L,b}^{\mathrm{EXACT}}(n_f=4)`$ (solid line) $`F_{L,b}^{\mathrm{CSN}}(n_f=5)`$, (dot-dashed line) together with the NLO charm density piece Term1 (dotted line) and the LO gluon density piece Term2 (dashed line), see text, for $`x=5\times 10^5`$ plotted as functions of $`Q^2`$.
The bottom quark structure function $`F_{L,b}^{\mathrm{EXACT}}(n_f=4)`$ (solid line) $`F_{L,b}^{\mathrm{CSN}}(n_f=5)`$, (dot-dashed line) together with the NLO charm density piece Term1 (dotted line) and the NLO gluon density piece Term2 (dashed line), see text, for $`x=5\times 10^5`$ plotted as functions of $`Q^2`$. |
warning/0002/astro-ph0002457.html | ar5iv | text | # Non-Uniform Reionization by Galaxies and its Effect on the Cosmic Microwave Background
## 1 Introduction
The Gunn-Peterson (GP) effect (?) strongly indicates that the smoothly distributed hydrogen in the intergalactic medium (IGM) is already highly ionized by $`z=5`$ (??). Barring the possibility of collisional reionization (e.g. Giroux & Shapiro 1994), the GP effect implies the presence of very luminous ionizing sources at high redshifts capable of producing enough Lyman continuum (Lyc) photons to cause photoionization of hydrogen by $`z>5`$. The two possible sources of these ionizing photons are QSOs and high mass stars.
Models in which QSOs dominate the production of ionizing photons may be able to meet the GP constraint (?). However, such models are strongly constrained by the observed drop in the abundance of bright QSOs above $`z3`$ (????). Furthermore, ? note that a model in which faint QSOs provide all the required ionizing luminosity can be ruled out on the basis of the number of faint QSOs seen in the HDF.
There is, however, growing evidence for the presence of bright galaxies at redshifts as high as $`z5`$ (?), and perhaps even higher (?). Thus the other natural candidate sources of ionizing photons are young, high mass stars forming in galaxies at redshifts greater than 5 (e.g. ???). ? note that at $`z3`$ stars in Lyman-break galaxies will emit more ionizing photons into the IGM than QSOs if more than 30% of such photons can escape from their host galaxy. Whilst such high escape fractions may not be realistic (e.g. local starbursts show escape fractions of only a few percent, ?), this does demonstrate that high-redshift galaxies could provide a significant contribution to (or perhaps even dominate) the production of ionizing photons. In this work we will restrict our attention to ionizing photons produced by stars, deferring consideration of the QSO contribution to a later paper.
According to the hierarchical structure formation scenario (e.g. ?) perturbations in the gravitationally dominant and dissipationless dark matter grow, by gravitational instability, into virialised clumps, or halos. Galaxies, and later stars, then form by the cooling and condensation of gas inside these halos (e.g. ??). Dark matter halos continually grow by merging with other halos (e.g. ??). In the context of this hierarchical scenario, we present a realistic scheme for studying the reionization of the universe by ionizing photons emitted from massive stars. We focus on the photoionization of the hydrogen component of the IGM. To predict the time dependent luminosity in Lyc photons we use a semi-analytic model of galaxy formation (e.g. ???). In particular, we use the semi-analytic model of ?, modified to take into account Compton cooling by cosmic microwave background (CMB) photons, to model the properties of galaxies living in dark matter halos spanning a wide range of masses.
We then estimate the fraction of the ionizing photons which manage to escape each galaxy, and therefore contribute to the photoionization of the intergalactic Hi. The fraction of ionizing photons escaping is determined on a galaxy-by-galaxy basis, using physically motivated models. Assuming spherical symmetry, we follow the propagation of the ionization front around each halo to compute the filling factor of intergalactic Hii regions, including the effects of clumping in the IGM. Finally, using several alternative models for the spatial distribution of ionized regions within a high resolution N-body simulation of the dark matter distribution, we estimate the anisotropies imprinted on the CMB by the patchy reionization process, due to the correlations in the ionized gas distribution and velocities (??). In previous models many simplifications were made in computing both the spatial distribution of ionized regions and the two-point correlations of gas density and velocity in those regions (??????). Our calculations represent a significant improvement over these models as we are able to calculate the two-point correlations between gas density and velocity in ionized regions directly from an N-body simulation.
The rest of this paper is arranged as follows. In §2 we outline the features of the semi-analytic model relevant to galaxy formation at high redshifts. In §3 we describe how we calculate the fraction of ionizing photons escaping from galaxies, and observational constraints on the ionizing luminosities and escape fraction at low and high redshift from $`\mathrm{H}\alpha `$ luminosities and Hi masses and column densities. In §4 we describe how we calculate the filling factor of photoionized gas in the IGM, including the effects of clumping of this gas. We then present our predictions for reionization, including the effects on the reionization redshift of using different assumptions about escape fractions and clumping factors.In §5 we examine the robustness of our results to changes in the other parameters of the semi-analytic galaxy formation model. In §6 we describe how the semi-analytic models are combined with N-body simulations to calculate the spatial distribution of the photoionized IGM. We then calculate the spectrum of anisotropies introduced into the CMB by this ionized gas. Finally, in §7 we summarize our results and examine their consequences.
## 2 The semi-analytic model of galaxy formation
To determine the luminosity in ionizing Lyc photons produced by the galaxy population, we use the semi-analytic model of galaxy formation developed by ?. This model predicts the properties of galaxies residing within dark matter halos of different masses. This is achieved by relating, in a self-consistent way, the physical processes of gas cooling, star formation, and supernovae feedback to a halo’s merger history, which is calculated using the extended Press-Schechter theory. The parameters of this model are constrained by a set of observations of galaxies in the local Universe, including the B and K-band luminosity functions, the I-band Tully-Fisher relation, the mixture of morphological types and the distribution of disk scale lengths (see ? and references therein for a thorough discussion of the observational constraints). Once the model has been constrained in this way it is able to make predictions concerning the clustering of galaxies (?) and the properties of galaxies at higher redshifts. For reference, the parameters of our standard model are given in Table 1. Definitions of the semi-analytic model parameters can be found in ?.
As we are employing the semi-analytic model at much higher redshifts than we have previously attempted, we will investigate the effects on our results of changing key model parameters. Of particular interest will be the prescription for feedback from supernovae and stellar winds. The model assumes that a mass $`\beta \mathrm{\Delta }M`$ of gas is reheated by supernovae and ejected from the disk for each mass $`\mathrm{\Delta }M`$ of stars formed. The quantity $`\beta `$ is allowed to be a function of the galaxy properties, and is parameterised as
$$\beta =(V_{\mathrm{d}isk}/V_{\mathrm{h}ot})^{\alpha _{\mathrm{h}ot}},$$
(1)
where $`V_{\mathrm{d}isk}`$ is the circular velocity of the galaxy disk and $`V_{\mathrm{h}ot}`$ and $`\alpha _{\mathrm{h}ot}`$ are adjustable parameters of the model. ? show that $`\alpha _{\mathrm{h}ot}`$ and $`V_{\mathrm{hot}}`$ are well constrained by the shape of the B-band luminosity function and the Tully-Fisher relation at $`z=0`$. However, since there is very little time available for star formation at the high redshifts which we consider, it is possible that these parameters, or indeed the form of the parameterisation in eqn. (1), could be changed at high redshift without significantly affecting the model predictions at $`z=0`$. In §5 we will therefore experiment with different values of these parameters and will also consider a modified functional form for $`\beta `$.
Two other key model inputs are the baryon density parameter, $`\mathrm{\Omega }_\mathrm{b}`$, and the stellar initial mass function (IMF). The value of $`\mathrm{\Omega }_\mathrm{b}`$ determines cooling rates (and so star formation rates) in our model halos. The shape of the IMF determines the number of high mass stars which produce the ionizing photons. For $`\mathrm{\Omega }_\mathrm{b}`$ our standard value is 0.02, which is consistent with the estimate ? and which allows a good match to the bright end of the observed B-band luminosity function. We will also consider an alternative value of $`\mathrm{\Omega }_\mathrm{b}=0.04`$, which is in better agreement with estimates from the D/H ratio in QSO absorption line systems (??). For the IMF, we adopt as our standard choice the IMF of ?, which is close to the “best” IMF proposed by ? on the basis of observations in the Solar neighbourhood and in nearby galaxies. We consider the effects of changing both $`\mathrm{\Omega }_\mathrm{b}`$ and the IMF in §5.
### 2.1 Gas Cooling
The standard semi-analytic model of ? allows hot halo gas to cool only via collisional radiative processes. At high redshift, Compton cooling due to free electrons in the hot plasma scattering off CMB photons becomes important. The Compton cooling timescale is given by (?)
$$t_{\mathrm{C}ompton}=\frac{1161.3(1+x_{\mathrm{h}ot}^1)}{(1+z)^4(1T_0^{\mathrm{C}MB}(1+z)/T_\mathrm{e})}\text{Gyr},$$
(2)
where $`x_{\mathrm{h}ot}`$ is the ionized fraction of the hot halo gas, $`T_0^{\mathrm{C}MB}`$ is the temperature of the CMB at the present day and $`T_\mathrm{e}`$ is the temperature of electrons in the hot halo gas, which we set equal to the virial temperature of the halo. At high redshifts this cooling time becomes shorter than the Hubble time and so Compton cooling may be effective at these redshifts. To implement eqn. (2) in the semi-analytic model, we assume that the shock-heated halo gas is in collisional ionization equilibrium, and use values for $`x_{\mathrm{h}ot}`$ which we interpolate from the tabulated values given by ?. In halos with virial temperatures less than around $`10^4`$K collisional ionization is ineffective, and so the ionized fraction in the halo gas will equal the residual ionization fraction left over from recombination. However, since this fraction is small, we will simply assume in this paper that cooling in halos below $`10^4`$K is negligible.
It should be noted that, unlike the radiative cooling time, the Compton cooling time is independent of the gas density and depends only very weakly on the gas temperature. Whereas with collisional radiative cooling a cooling radius, within which the cooling time is less than the age of the halo, propagates through the halo as more and more gas cools, with Compton cooling the entire halo cools at the same rate. The amount of gas which can reach the centre of the halo is then controlled by the free-fall timescale in the halo.
Including Compton cooling in our model turns out to make little difference to the results. For example, the total mass of stars formed in the universe as a function of redshift differs by less than $`5\%`$ for $`z<20`$ between models with and without Compton cooling. At higher redshifts the differences can become as large as 30% for small intervals of redshift. For example, if a massive halo cools via Compton cooling it will rapidly produce many stars — without Compton cooling it will still form these stars, but not until slightly later when collisional radiative cooling takes effect. However, the mass of stars formed at these high redshifts is tiny, and so any differences become entirely negligible at lower redshifts when many more stars have formed. Although at high redshift the Compton cooling time is shorter than the age of the Universe, halos are merging at a high rate and so their gas is being repeatedly shock-heated by successive mergers which, we assume, heat the gas to the virial temperature of the halo. We find that for the majority of halos the time between successive major mergers (defined as the time for a doubling in mass of the halo) is less than the Compton cooling time at the redshifts considered here. Therefore, Compton cooling will be ineffective in these halos. In the few cases where the halo does survive long enough that Compton cooling could be important, we find that the collisional radiative cooling time at the virial radius of the halo is often shorter than the Compton cooling time, in which case all of the gas in the halo will cool whether or not we include the effects of Compton cooling. Nevertheless, Compton cooling is included in all models considered.
In this work we ignore cooling due to molecular hydrogen (H<sub>2</sub>). Although molecular hydrogen allows cooling to occur in gas below $`10^4`$K, it is easily dissociated by photons from stars that form from the cooling gas. Previous studies that have included cooling due to H<sub>2</sub> typically find that it is completely dissociated at very high redshifts. For example, ? find that molecular hydrogen is fully dissociated by $`z25`$. Objects formed by H<sub>2</sub> cooling are therefore not expected to contribute significantly to the reionization of the IGM.
### 2.2 Fraction of Gas in the IGM
At any given redshift, some fraction of the gas in the Universe will have become collisionally ionized in dark matter halos and some fraction will have cooled to become part of a galaxy. Within the context of our semi-analytic model, we define the IGM as all gas which has *not* been collisionally ionized inside dark matter halos and which has *not* become part of a galaxy (note that we are here only interested in ionization of hydrogen). It is this gas which must be photoionized if the Gunn-Peterson constraint is to be satisfied. The fraction of the total baryon content of the universe which is in the IGM, $`f_{\mathrm{IGM}}`$, can be estimated by integrating over the mass function of dark matter halos, as follows
$$f_{\mathrm{IGM}}(z)=1\left[_0^{\mathrm{}}M_{\mathrm{gas}}x_\mathrm{H}\frac{\mathrm{d}n}{\mathrm{d}M_{\mathrm{halo}}}\frac{\mathrm{d}M_{\mathrm{halo}}}{\mathrm{\Omega }_\mathrm{b}\rho _\mathrm{c}}\right]f_{\mathrm{galaxy}}(z),$$
(3)
where $`M_{\mathrm{gas}}(M_{\mathrm{halo}},z)`$ is the mean mass of diffuse gas in halos of mass $`M_{\mathrm{halo}}`$, $`x_\mathrm{H}(M_{\mathrm{halo}},z)`$ is the fraction of hydrogen which is collisionally ionized at the halo virial temperature (which we take from the calculations of ?), $`\mathrm{d}n/\mathrm{d}M_{\mathrm{halo}}(M_{\mathrm{halo}},z)`$ is the comoving number density of halos (which we approximate by the Press-Schechter mass function), $`\rho _\mathrm{c}`$ is the critical density of the Universe at $`z=0`$, and $`f_{\mathrm{galaxy}}`$ is the fraction of the total baryonic mass in the Universe which has been incorporated into galaxies. The quantities $`M_{\mathrm{gas}}`$ and $`f_{\mathrm{galaxy}}`$ can be readily calculated from our model of galaxy formation. Fig. 1 shows the evolution of $`f_{\mathrm{IGM}}`$ with redshift.
### 2.3 Observational Constraints
The semi-analytic model provides the spectral energy distribution (SED) of each galaxy, from which we can determine the ionizing luminosity of that galaxy. Summing the contributions from all galaxies in a given halo yields the total ionizing luminosity produced in that halo. ?, using the model of ?, demonstrated that a higher reionization redshift could be obtained if zero-metallicity stars were responsible for reionization, as these produce a greater number of ionizing photons than low (i.e. $`10^4`$) metallicity stars. In our model the very first stars have zero metallicity, but as we include chemical evolution only a very small fraction of stars have metallicities below $`10^4`$. This is consistent with the results of ? who argue that the epoch of metal-free star formation must end before $`z=3`$, as the enhanced emission shortwards of 228Å from such stars is inconsistent with observations of Heii opacity in the IGM at that redshift. Therefore we cannot appeal to such zero-metallicity stars to increase the redshift of reionization in our model.
As a result of absorption by neutral hydrogen close to the emitting stars and extinction caused by dust, only a small fraction of the ionizing radiation emitted by the stars escapes from each galaxy (???). We therefore estimate, within the context of the semi-analytic model, the fraction, $`f_{\mathrm{esc}}`$, of ionizing photons which escape the galaxy to become available for the photoionization of the Hi in the IGM. The calculation of $`f_{\mathrm{esc}}`$ is discussed in §3.
In Fig. 2 we show the redshift evolution of the comoving number density of galaxies with ionizing luminosity $`\dot{n}_{\mathrm{i}on}`$ larger than $`10^2`$, $`10^3`$ and $`10^4`$ in units of $`10^{50}`$ photons per second. These are the unattenuated luminosities produced by massive stars in the galaxies. The abundances of sources of given luminosity rises sharply up to $`z=24`$ (the exact position of the peak depending on luminosity) as more and more dark matter halos form that are capable of hosting bright galaxies. After $`z=24`$ abundances quickly drop towards $`z=0`$ as the amount of gas available for star formation declines.
The escape fractions in our model will be determined by the mass and radial scale length of the Hi gas in galactic disks. It is therefore important to test that our model produces galaxies with reasonable distributions of Hi mass and disk scale length. ? have shown that our semi-analytic model produces distributions of I-band disk scale lengths in good agreement with the $`z=0`$ data of ?. In Fig. 3 we compare our model with observations of damped Lyman-$`\alpha `$ systems (DLAS) over a range of redshifts and with the Hi mass function at $`z=0.0`$. Under the assumption that DLAS are caused by neutral gas in galactic disks, we compute the DLAS column density distribution in our model, $`f_{\mathrm{DLAS}}`$, defined such that $`f_{\mathrm{DLAS}}(N_{\mathrm{H}\mathrm{i}},t)\mathrm{d}N_{\mathrm{H}\mathrm{i}}\mathrm{d}X`$ is the mean number of DLAS at cosmic time $`t`$ with column densities in the range $`N_{\mathrm{H}\mathrm{i}}`$ to $`N_{\mathrm{H}\mathrm{i}}+\mathrm{d}N_{\mathrm{H}\mathrm{i}}`$ and absorption distance $`X(z)=\frac{2}{3}\left[(1+z)^{3/2}1\right]`$ in the interval $`\mathrm{d}X`$ along a line of sight (?). Our model is in reasonable agreement with the distribution of DLAS column densities observed by ?, indicating that both the mass of Hi and its radial scale length in our model galaxies are realistic. The $`z=0`$ Hi mass function from our model is also in reasonable agreement with the data of ?, although it does overpredict the abundance of low Hi mass galaxies. Our model predictions assume that all of the hydrogen in galactic disks is in the form of Hi. In practice, some of the hydrogen in disks will be in the form of molecules (H<sub>2</sub>) or ionized gas (Hii), so this over-estimates the Hi masses and column densities.
A significant contribution to the ionizing luminosity comes from very low mass halos. We therefore ensure that we resolve all halos which have a virial temperature $`10^4`$K up to $`z=50`$, i.e. all halos down to a mass of $`5\times 10^6h^1M_{}`$. Below this temperature cooling becomes inefficient (since we are ignoring cooling by molecular hydrogen, and the Compton cooling from the residual free electrons left over after recombination) and so galaxy formation ceases.
The requirement that $`5\times 10^6h^1M_{}`$ halos be resolved sets an upper limit on the mass of halo that we can simulate due to computer memory limits, since the lower the mass of halo that is resolved, the more progenitors a halo of given mass will have. At $`z=0`$, the most massive halos that we are able to simulate make a significant contribution to the total filling factor and ionizing luminosity. However, for $`z2`$ the most massive halos simulated contribute only 1% of the total number of escaping ionizing photons, and this fraction drops extremely rapidly as we look to even higher redshifts. Therefore, at the high redshifts ($`z>3`$) we will be interested in, ignoring higher mass halos makes no significant difference to our results.
We note that once any halo has begun to ionize the surrounding IGM, it could potentially influence the process of galaxy formation in nearby halos. Ionizing photons from the first halo will act to heat the gas in nearby halos, thereby reducing the effective cooling rate (??). Since prior to full reionization each halo will see only the flux of ionizing photons from nearby sources, a detailed accounting of this radiative feedback requires a treatment of the radiative transfer of the ionizing radiation through the IGM. This is beyond the scope of the present work. Such radiative feedback is expected to be very efficient at dissociating molecular hydrogen, with ? finding that H<sub>2</sub> is completely dissociated by $`z25`$. Radiative feedback will also inhibit galaxy formation both by reducing the amount of gas that accretes into low mass halos (?) and by reducing the cooling rate of gas within halos. ? show that radiative feedback may be effective in inhibiting galaxy formation in halos with circular velocities of 50 km/s or less. In our model, the ionizing luminosity becomes dominated by galaxies in halos with circular velocities greater than 50 km/s at redshifts below $`z10`$. At higher redshifts we may therefore be overestimating the total ionizing luminosity produced by galaxies, but this should not significantly affect the reionization redshift.
## 3 The escape fraction of ionizing photons
### 3.1 Global constraints at low redshift
Gas and dust inside galaxies can readily absorb ionizing photons and re-emit the energy at longer wavelengths. Therefore the amount and distribution of these components are the main factors that determine $`f_{\mathrm{esc}}`$. The model of galaxy formation explicitly provides the mass and metallicity of cold gas present in each galaxy disk and the half-mass radius of that disk, all as functions of time. The mass of dust is assumed to be proportional to the mass of cold gas and to its metallicity. We split the escape fraction into contributions from gas, $`f_{\mathrm{esc},\mathrm{gas}}`$, and dust, $`f_{\mathrm{esc},\mathrm{dust}}`$, such that the total escaping fraction is given by $`f_{\mathrm{esc}}=f_{\mathrm{esc},\mathrm{gas}}f_{\mathrm{esc},\mathrm{dust}}`$.
? describe in detail how the effects of dust are included in their model of galaxy formation. This modelling, which uses the calculations of ?, is much more realistic than has been previously included in semi-analytic galaxy formation models, as it includes a fully 3D (though axisymmetric) dust distribution, and the dust optical depths are calculated for each galaxy individually. In this model, stars are assumed to be distributed in a bulge and in an exponential disk with a vertical scale height equal to $`0.0875`$ times the radial scale length (this ratio was adopted by Ferrara et al. to match the observed values for the old disk population of galaxies like the Milky Way). The dust is assumed to be distributed in the same way as the disk stars. The models give the attenuation of the ionizing radiation as a function of the inclination angle at which a galaxy is viewed, and we average this over angle to find the mean dust extinction for each galaxy. The dust attenuations do not include the effects of clumping of the stars or dust, and also assume that the ionizing stars have the same vertical distribution as the dust. With these two caveats in mind, the dust extinctions we apply should only be considered as approximate.
Some of the emitted Lyc photons are absorbed by neutral hydrogen close to the emitting star, thereby causing H$`\alpha `$ line emission from the galaxy. Therefore, the H$`\alpha `$ luminosity function is sensitive to the fraction, $`f_{\mathrm{esc},\mathrm{gas}}`$, of the ionizing photons which manage to escape through the gas. We will require our models to reproduce the observed H$`\alpha `$ luminosity function and luminosity density.
In Figs. 4 and 5 we compare the H$`\alpha `$ emission line properties of galaxies in our model with observational data at low redshift. The observed values are already corrected for dust extinction, so we compare them with the theoretical values before dust attenuation. In order to calculate these properties accurately, we simulate halos of mass up to and including $`10^{15}h^1M_{}`$. In calculating the H$`\alpha `$ line luminosity of each galaxy, we assume that a fraction $`1f_{\mathrm{e}sc,gas}`$ of the Lyc photons are absorbed by hydrogen atoms, producing H$`\alpha `$ photons according to case B recombination. The remaining Lyc photons escape, after being further attenuated by dust. The figures show results for $`f_{\mathrm{e}sc,gas}=0`$, 0.05 and 0.2, which roughly brackets the likely range of values for typical disk galaxies at the present day, as we discuss below. Both the predicted H$`\alpha `$ luminosity function and luminosity density are in reasonable agreement with the observations, demonstrating that our models produce galaxies with realistic total ionizing luminosities (before attenuation by gas and dust). In principle, these observational comparisons provide a constraint on the value of $`f_{\mathrm{esc},\mathrm{gas}}`$, if the other parameters in the semi-analytical model are assumed to be known. However, in practice it is not possible to reliably distinguish between $`f_{\mathrm{esc},\mathrm{gas}}=0.2`$ and $`f_{\mathrm{esc},\mathrm{gas}}=0`$ or $`0.05`$, given the uncertainties in the observational data. The observational results depend on the dust correction factors applied, and there is also some uncertainty in the ionizing luminosities predicted by stellar population synthesis models for a given IMF. With these caveats in mind it would seem that mean escape fractions anywhere between zero and 20-30% are acceptable. Observations of starburst galaxies in the nearby universe suggest that the escape fraction is actually less than $`3\%`$ for such galaxies (?), but starbursts are known to have very high column densities of gas and dust, and so the escape fraction in normal galaxies can probably be significantly higher (e.g. ?).
### 3.2 The dependence of $`𝐟_{𝐞sc,gas}`$ on redshift and on halo mass
So far, we have assumed that $`f_{\mathrm{esc},\mathrm{gas}}`$ is a global constant, varying neither with galaxy properties nor redshift. The details of the physical processes which determine $`f_{\mathrm{esc},\mathrm{gas}}`$ are uncertain, but a constant $`f_{\mathrm{esc},\mathrm{gas}}`$ seems unrealistic, as the properties of the emitting galaxies depend strongly upon both redshift and the mass of the halos in which they live. Given the complexity of this problem, here we merely aim at establishing the general trend of how $`f_{\mathrm{esc},\mathrm{gas}}`$ may vary with halo mass and redshift.
We will consider three models for $`f_{\mathrm{esc},\mathrm{gas}}`$. In the first model, $`f_{\mathrm{esc},\mathrm{gas}}`$ is assumed to be a universal constant (this will be referred to as the “fixed model”). In the second and third models $`f_{\mathrm{esc},\mathrm{gas}}`$ is evaluated for each galaxy, based on its physical properties. These two models are described next.
Our first physical $`f_{\mathrm{esc},\mathrm{gas}}`$ model is based on the approach of ?, hereafter DS94 who derived an analytic expression for $`f_{\mathrm{esc},\mathrm{gas}}`$. In their model, Lyc photons are emitted by OB associations in a galactic disk and escape by ionizing “Hii chimneys” in the Hi layer. The fraction of photons escaping a disk of given size and gas content can then be calculated. Whilst the original DS94 model assumes that OB associations all lie in the mid-plane of the galaxy disk, we have also considered the case where OB associations are distributed vertically like the gas in the disk.
The Hii chimney model of DS94 does not include the effects of finite lifetimes of the OB associations, or of dynamical evolution of the gas distribution around an OB association due to energy input by stellar winds and supernova from the OB association itself. ? have calculated the escape of ionizing photons through a dynamically evolving superbubble, which is driven by an OB association at its centre. They find that the resulting escape fractions are slightly lower than those obtained from the DS94 model (since the superbubble shell is able to effectively trap radiation). Numerical solutions of the radiative transfer equations in disk galaxies give results in excellent agreement with the Strömgren sphere approach of DS94 for OB associations at the bright end of the luminosity function, but give somewhat lower escape fractions for the faintest OB associations, the two approaches differing by around 25% for a single OB star (?).
Our second physical model for $`f_{\mathrm{esc},\mathrm{gas}}`$ is based on ?, hereafter DSGN98. In this case, the ionizing stars are assumed to be uniformly mixed with the gas in the galaxy, and the gas is assumed to remain neutral. DSGN98 give an approximate analytic expression for the escape fraction in this case, but we have instead calculated the escape fraction exactly by numerical integration, for a specific choice for the gas density profile.
We give details of the calculation of $`f_{\mathrm{esc},\mathrm{gas}}`$ in the DS94 and DSGN98 models in Appendix A. Both models contain one free parameter, $`h_\mathrm{z}/r_{\mathrm{disk}}`$, the ratio of disk scale height to radial scale length. We will consider the effects of varying this parameter in §5. For starbursts, we calculate the escape fraction based on a simple spherical geometry, as is also described in Appendix A. The contribution to the total ionizing luminosity from bursts of star formation is small ($`<8\%`$) at all redshifts.
To summarize, we will show results from three models for $`f_{\mathrm{esc},\mathrm{gas}}`$ as standard. These are: a model in which $`f_{\mathrm{esc},\mathrm{gas}}`$ is held constant at 0.1; the DS94 model with OB associations in the disk midplane; and the DSGN98 model using our exact calculation of the escaping fraction. We consider the DS94 model to be the most realistic of our three models for $`f_{\mathrm{esc},\mathrm{gas}}`$, but also present results from the other models for comparison.
In Fig. 6 we show the variation of $`f_{\mathrm{esc}}`$ with dark halo mass at $`z=0`$ for the three models. The thin and thick lines show the escape fraction respectively with and without attenuation by dust. When a halo contains more than one galaxy, we plot the mean $`f_{\mathrm{esc}}`$ weighted by ionizing luminosity. At a given halo mass, halos with the lowest $`f_{\mathrm{esc}}`$ tend to have the highest ionizing luminosities, as both the star formation rate and attenuation of photons are increased in galaxies with large gas contents. The three models all show a trend for decreasing escape fraction with increasing halo mass up to $`M_{\mathrm{halo}}10^{12}h^1M_{}`$. For the fixed gas escape fraction model, the variation in $`f_{\mathrm{esc}}`$ is due entirely to the effects of dust, which can therefore be seen to be negligible in halos less massive than $`10^{10}h^1M_{}`$. This decrease in $`f_{\mathrm{esc}}`$ due to dust is enhanced in the other two models by the variation in $`f_{\mathrm{esc},\mathrm{gas}}`$, which also declines with increasing halo mass. For halos more massive than $`10^{12}h^1M_{}`$, the escape fractions rise somewhat for the variable $`f_{\mathrm{esc},\mathrm{gas}}`$ models. Note that the DSGN98 model predicts a much smaller escape fraction than the DS94 model at all masses.
In Fig. 7 we plot the variation in the (ionizing luminosity-weighted) mean disk scale length and cold gas mass for galaxies in our model as a function of halo mass. Evidently, the decline in $`f_{\mathrm{esc}}`$ with increasing halo mass below $`10^{12}h^1M_{}`$ seen in Fig. 6 is due mainly to the greater masses of gas found in galaxies in these halos. This rapid change in the mass of gas present is due to the effects of feedback, which efficiently ejects gas from galaxies in low mass halos. Above $`10^{12}h^1M_{}`$, the mass of cold gas in galaxies levels off and then begins to decline as cooling becomes inefficient in more massive halos. This results in an escape fraction increasing with halo mass for the most massive halos simulated. Although galaxy sizes increase with increasing halo mass, thereby reducing gas densities somewhat, this effect is not strong enough to offset the increased cold gas mass in these galaxies.
We find that the DS94 and DSGN98 models applied to our galaxies predict escape fractions (including the effects of dust) for halos of mass $`10^{11}h^1M_{}`$ at $`z=0`$ of $`20\%`$ and $`0.2\%`$ respectively. The mean DS94 and DSGN98 luminosity-weighted escape fractions for galaxies at $`z=0`$ are lower, being $`6\%`$ and $`0.1\%`$ respectively. However, we expect some variation in these values with redshift due to the evolution of the galaxy population. In fact, we find a rapid decline in both cold gas content and galaxy disk size with increasing redshift. In Fig. 8 we show the evolution of the (ionizing luminosity-weighted) mean $`f_{\mathrm{esc}}`$ between redshifts 0 and 45. All models show an initial rapid decline in $`f_{\mathrm{esc}}`$ with increasing $`z`$. After this, in the constant $`f_{\mathrm{esc},\mathrm{gas}}`$ model, the mean escape fraction increases with redshift since the dust content of galaxies was lower in the past. The DS94 model shows a very gradually rising escape fraction, whilst the DSGN98 model has a more rapid decline.
In our model, the contribution of stellar sources to the UV background is dominated by galaxies at low redshifts ($`z<1`$). We find that immediately shortwards of 912Å our DS94 model predicts a background due to stellar sources which is very close to that expected from QSOs (?), after including the effects of attenuation by the intervening IGM (?). At shorter wavelengths the QSO contribution soon becomes dominant. Thus, at $`z=0`$ the combined background due to stars (from our DS94 model) and QSOs (from ?) is $`J_{\text{912Å}}4\times 10^{23}`$ ergs/s/cm<sup>2</sup>/Hz/ster. This is consistent with the upper limit of $`J_{\text{912Å}}=8\times 10^{23}`$ ergs/s/cm<sup>2</sup>/Hz/ster found by ?, who searched for H$`\alpha `$ emission from intergalactic HI clouds. The contribution of galaxies to the local ionizing background has also been estimated by ?, based on the luminosity function of galaxies observed in the Canada-France Redshift Survey. They estimated the galactic contribution as $`J_{\text{912Å}}5\times 10^{23}`$ ergs/s/cm<sup>2</sup>/Hz/ster at $`z=0`$, assuming an escape fraction of $`f_{\mathrm{esc}}=0.15`$. If we assume the same $`f_{\mathrm{esc}}`$ in our model, we obtain $`J_{\text{912Å}}=5.2\times 10^{23}`$ ergs/s/cm<sup>2</sup>/Hz/ster, in excellent agreement with their result.
## 4 The filling factor and the evolution of the ionization fronts
We define the filling factor, $`F_{\mathrm{fill}}`$, as the fraction of hydrogen in the IGM (as defined in §2.2) which has been ionized. This is the natural quantity which serves as an indication of the amount of reionization in the IGM. We calculate the growth of the ionized region around each halo, using the ionizing luminosities predicted by the semi-analytic model, and then sum over all halos to find $`F_{\mathrm{fill}}`$. We make two simplifying assumptions: (1) the radiation from each halo is emitted isotropically, and (2) the distribution of hydrogen is uniform on the scale of the ionization front and larger (but with small-scale clumping). It follows that each halo by itself would produce a spherical ionization front.
The mass of hydrogen ionized within the ionization front, $`M`$, in spherical symmetry is given by (??)
$$\frac{1}{\mathrm{m}_\mathrm{H}}\frac{\mathrm{d}M}{\mathrm{d}t}=S(t)\alpha _\mathrm{H}^{(2)}a^3f_{\mathrm{clump}}n_\mathrm{H}\frac{M}{\mathrm{m}_\mathrm{H}},$$
(4)
where $`n_\mathrm{H}`$ is the comoving mean number density of hydrogen atoms (total, Hi and Hii) in the IGM, $`a(t)`$ is the scale factor of the universe normalized to unity at $`z=0`$, $`t`$ is time and $`S(t)`$ is the rate at which ionizing photons are being emitted. The factor $`f_{\mathrm{clump}}(t)`$, defined by
$$f_{\mathrm{clump}}=\rho _{\mathrm{IGM}}^2/\overline{\rho }_{\mathrm{IGM}}^2,$$
(5)
is the clumping factor for the ionized gas in the IGM (here $`\rho _{\mathrm{IGM}}`$ is the density of IGM gas at any point and $`\overline{\rho }_{\mathrm{IGM}}`$ is the mean density of the IGM). Small-scale clumpiness causes the total recombination rate to be larger than for a uniform medium of the same mean density.
The value of $`f_{\mathrm{c}lump}(t)`$ for the ionized gas in the IGM is complicated to calculate analytically. We remind the reader that in our picture, the IGM consists of all gas which has *not* been collisionally ionized in halos nor become part of a galaxy. For a uniform IGM $`f_{\mathrm{c}lump}=1`$ by definition. If low density regions of the IGM are ionized before high density regions, as suggested by ?, then this would be similar to having $`f_{\mathrm{c}lump}<1`$ in eqn. (4), but for most purposes, $`f_{\mathrm{c}lump}=1`$ can be considered as an approximate lower bound. We make two different estimates of the possible effects of clumping.
For our first estimate, which we call $`f_{\mathrm{clump}}^{(\mathrm{variance})}`$, we assume that the photoionized gas basically traces the dark matter, except that gas pressure prevents it from falling into dark matter halos with virial temperatures smaller than $`10^4`$K (the approximate temperature of the photo-ionized gas). Thus, we calculate the clumping factor as $`f_{\mathrm{clump}}^{(\mathrm{variance})}=(1+\sigma ^2)`$, where $`\sigma ^2`$ is the variance of the dark matter density field in spheres of radius equal to the virial radius of a $`10^4`$K halo. $`\sigma ^2`$ is calculated from the non-linear dark matter power spectrum, estimated using the procedure of ?, and smoothed using a top-hat filter in real space.
For our second estimate, which we call $`f_{\mathrm{clump}}^{(\mathrm{halos})}`$, we include the effects of collisional ionization in halos and of removal of gas by cooling into galaxies in a way consistent with our definition of $`f_{\mathrm{IGM}}`$ given in equation (3). The diffuse gas in halos with virial temperatures above $`10^4`$K is assumed to have the density profile of an isothermal sphere with a constant density core. The gas originally associated with smaller halos is assumed to be pushed out of these halos by gas pressure following photoionization, and to be in a uniform density component occupying the remaining volume. As shown in Appendix B, the clumping factor is then
$`f_{\mathrm{clump}}^{(\mathrm{halos})}`$ $`=`$ $`{\displaystyle \frac{f_{\mathrm{m},\mathrm{smooth}}^2}{f_{\mathrm{v},\mathrm{smooth}}f_{\mathrm{IGM}}^2}}+{\displaystyle \frac{f_{\mathrm{int}}\mathrm{\Delta }_{\mathrm{vir}}}{f_{\mathrm{IGM}}^2}}{\displaystyle _{M_\mathrm{J}}^{\mathrm{}}}(1f_{\mathrm{gal}})^2`$ (6)
$`\times (1x_\mathrm{H})^2{\displaystyle \frac{M_{\mathrm{halo}}}{\rho _\mathrm{c}\mathrm{\Omega }_0}}{\displaystyle \frac{\mathrm{d}n}{\mathrm{d}M_{\mathrm{halo}}}}\mathrm{d}M_{\mathrm{halo}},`$
where $`M_\mathrm{J}`$ is the mass of a halo which just retains reionized gas, $`f_{\mathrm{m},\mathrm{smooth}}`$ is the fraction of the total baryonic mass in the uniform component, and $`f_{\mathrm{v},\mathrm{smooth}}`$ is the fraction of the volume of the universe occupied by this gas. Here, $`f_{\mathrm{gal}}`$ is the fraction of the baryonic mass in a halo in the form of galaxies, $`x_\mathrm{H}`$ is the fraction of hydrogen in the diffuse halo gas which is collisionally ionized (as in eqn. 3), and $``$ indicates an average over all halos of mass $`M_{\mathrm{halo}}`$. The factor $`f_{\mathrm{int}}`$ is a parameter depending only on the ratio of the size of the core in the gas density profile to the halo virial radius. For a core radius equal to one-tenth of the virial radius $`f_{\mathrm{int}}=3.14`$ (see Appendix B). This estimate of the clumping factor ignores the possibility of gas in the centres of halos (but not part of a galaxy) becoming self-shielded from the ionizing radiation. Such gas would not become photoionized, and so would not contribute to the recombination rate, resulting in $`f_{\mathrm{clump}}`$ being lower than estimated here. A detailed treatment of the ionization and temperature structure of gas inside halos is beyond the scope of this work.
Of course, before reionization, when the gas is typically much cooler than $`10^4`$K, gas will fall into dark matter halos with virial temperatures below $`10^4`$K. If all gas were in halos, then we would find
$`f_{\mathrm{clump}}`$ $``$ $`{\displaystyle \frac{f_{\mathrm{int}}\mathrm{\Delta }_{\mathrm{vir}}}{f_{\mathrm{IGM}}^2}}{\displaystyle _0^{\mathrm{}}}(1f_{\mathrm{gal}})^2(1x_\mathrm{H})^2{\displaystyle \frac{M_{\mathrm{halo}}}{\rho _\mathrm{c}\mathrm{\Omega }_0}}`$ (7)
$`\times {\displaystyle \frac{\mathrm{d}n}{\mathrm{d}M_{\mathrm{halo}}}}\mathrm{d}M_{\mathrm{halo}}`$
Once reionized, some of the gas in these small halos will flow back out of the halo as the gravitational potential is no longer deep enough to confine the gas.
We consider $`f_{\mathrm{clump}}^{(\mathrm{halos})}`$ as our best estimate of the IGM clumping factor, at least for lower redshifts, when a significant fraction of the gas is in halos with $`M>M_J`$, while $`f_{\mathrm{clump}}^{(\mathrm{variance})}`$ is more in the nature of an upper limit.
The clumping factors $`f_{\mathrm{clump}}^{(\mathrm{halos})}`$ and $`f_{\mathrm{clump}}^{(\mathrm{variance})}`$ are plotted as functions of redshift in Fig. 9. They show fairly similar behaviour above $`z10`$. Below this redshift, $`f_{\mathrm{clump}}^{(\mathrm{variance})}`$ greatly exceeds $`f_{\mathrm{clump}}^{(\mathrm{halos})}`$, because it becomes dominated by gas in massive dark matter halos, which, on the other hand, contributes negligibly to $`f_{\mathrm{clump}}^{(\mathrm{halos})}`$ as it is collisionally ionized. We also plot estimates of the clumping factor from two other papers: ? calculated the clumping factor of the baryons which have been unable to cool (the quantity they call $`C_n`$) using their own analytical model. They obtain values of $`f_{\mathrm{clump}}`$ which are comparable to $`f_{\mathrm{clump}}^{(\mathrm{halos})}`$ at $`z<10`$, but are substantially larger at higher redshifts. ? performed hydrodynamical simulations in a cosmology similar to that which we consider, from which they measured $`f_{\mathrm{clump}}`$ directly. They calculated two clumping factors: one for all baryons in their simulation, $`f_{\mathrm{clump}}=f_{\mathrm{clump}}^{(\mathrm{GO}:\mathrm{bb})}`$, and the other for baryons in ionized regions only, $`f_{\mathrm{clump}}=f_{\mathrm{clump}}^{(\mathrm{GO}:\mathrm{H}_{\mathrm{II}})}`$, which is smaller. $`f_{\mathrm{clump}}^{(\mathrm{GO}:\mathrm{H}_{\mathrm{II}})}`$ is more relevant for our purposes, but may still overestimate the clumping of photo-ionized gas in the IGM, since it includes collisionally ionized gas in galaxy halos. These clumping factors are everywhere lower than $`f_{\mathrm{clump}}=f_{\mathrm{clump}}^{(\mathrm{variance})}`$. $`f_{\mathrm{clump}}^{(\mathrm{GO}:\mathrm{H}_{\mathrm{II}})}`$ is close to our estimate $`f_{\mathrm{clump}}^{(\mathrm{halos})}`$ at the highest and lowest redshifts, but smaller in the intermediate range.
### 4.1 Model results
In Fig. 10 we show the ionized filling factor of the IGM, $`F_{\mathrm{fill}}`$, as a function of redshift. Here we compute $`F_{\mathrm{fill}}`$ by summing the volumes of the Hii regions formed around each halo, weighted by the number of such halos per unit volume as given by the Press-Schechter theory. (Later we will use the halo mass function measured directly from an N-body simulation to calculate $`F_{\mathrm{fill}}`$ — see Fig. 14). $`F_{\mathrm{fill}}`$ will exceed 1 if more ionizing photons have been produced than are needed to completely reionize the universe. We show results for our three models for $`f_{\mathrm{esc},\mathrm{gas}}`$, and for three different assumptions about $`f_{\mathrm{clump}}`$.
If we ignore the effects of dust, we find that the model with a constant escape fraction of 10% reionizes the Universe by $`z=7.9`$ if $`f_{\mathrm{clump}}=1`$ but only by $`z=6.6`$ if $`f_{\mathrm{clump}}=f_{\mathrm{clump}}^{(\mathrm{halos})}`$. In order to reionize the Universe by $`z=5`$, escape fractions of 1.4% and 3.7% are needed for $`f_{\mathrm{clump}}=1`$ and $`f_{\mathrm{clump}}=f_{\mathrm{clump}}^{(\mathrm{halos})}`$ respectively. When we include the effects of dust, we find that gas escape fractions of 3.3% and 9.3% are needed to reionize by $`z=5`$ for these two cases.
If, instead of assuming a constant gas escape fraction, we use the more physically motivated DS94 model, we find reionization occurs at $`z=6.1`$ if $`f_{\mathrm{clump}}=1`$, but only at $`z=4.5`$ if $`f_{\mathrm{clump}}=f_{\mathrm{clump}}^{(\mathrm{halos})}`$ (both estimates including dust). In the latter case, the ionized filling factor at $`z=5`$ is only 76%. If we assume OB associations are distributed as the gas in the DS94 model (as opposed to lying in the disk mid-plane as in our standard model) then a filling factor of 117% (i.e. full reionization) is achieved by $`z=5`$. The DSGN98 model, which predicts much lower escape fractions than the DS94 model at all redshifts, is able to reionize only $`2`$% of the IGM by redshift 3, and full reionization never occurs even if $`f_{\mathrm{clump}}=1`$.
We note that in the DS94 model, approximately 90% of the photons required for ionization are produced at $`z<10`$. Thus our neglect of radiative feedback effects (which may reduce the number of ionizing photons produced at higher redshifts) is unlikely to seriously effect our determination of the reionization epoch.
As both the DS94 and DSGN98 models predict quite low escape fractions, we have also considered a much more extreme model which simply assumes that $`f_{\mathrm{esc}}=\beta /(1+\beta )`$, where $`\beta `$ is the feedback efficiency as defined by eqn. (1). This toy model, which we will refer to as the “holes scenario”, produces very high escaping fractions for galaxies with low circular speeds, and low escaping fractions for those with high circular speeds. A behaviour for $`f_{\mathrm{esc}}`$ of this general form might result if photons are able to escape through holes in the galaxy disk which have been created by supernovae. Since dust would also be expected to be swept out of these holes we do not include any dust absorption in this model. The holes scenario produces very different results compared to our two physical models for the escape fraction. In this model $`f_{\mathrm{esc}}1`$ for $`z>10`$, dropping to 45% by $`z=0`$. Not surprisingly therefore, this model succeeds in satisfying the Gunn-Peterson constraint, reionizing the Universe by $`z=11.7`$ if $`f_{\mathrm{clump}}=1`$ and by $`z=10.6`$ if $`f_{\mathrm{clump}}=f_{\mathrm{clump}}^{(\mathrm{halos})}`$. While this model is only a very crude attempt to consider a dynamically disturbed gas distribution in galaxy disks, it clearly demonstrates that such effects may be of great importance for studies of reionization.
We have also computed the filling factor in our model using the clumping factors calculated by ? and ? (as given in Fig. 9). Of course, this is not strictly self-consistent, as their clumping factors are calculated from their own models for galaxy formation and reionization, which differ from ours. Using either of these with the DS94 model gives a reionization redshift comparable to that obtained using $`f_{\mathrm{clump}}=f_{\mathrm{clump}}^{(\mathrm{halos})}`$: we find reionization at $`z=3.6`$ using the ? clumping factor, and $`z=4.9`$ using that of ?.
In Fig. 11 we show the total number of ionizing photons which have escaped into the IGM per unit comoving volume by redshift $`z`$, $`n_\gamma `$, divided by the total number of hydrogen nuclei in the IGM per unit comoving volume, $`n_\mathrm{H}`$ (which is $`f_{\mathrm{IGM}}`$ times the total number density of hydrogen nuclei). When this number reaches one, just enough photons have been emitted by galaxies to reionize the IGM completely if recombinations are unimportant. This criterion has been used previously to estimate when reionization may occur. Since our model includes the effects of recombinations in the IGM, we can judge how well this simpler criterion performs. If we ignore the effects of absorption by gas and dust on the number of ionizing photons escaping from galaxies, we find that, in this cosmology, our model achieves $`n_\gamma /n_\mathrm{H}=1`$ by $`z12`$. When we account for the effects of dust and gas in galaxies, we find that the redshift at which $`n_\gamma /n_\mathrm{H}=1`$ is significantly reduced, the exact value depending on the model for the escape fraction. With a fixed gas escape fraction of $`10\%`$, $`n_\gamma /n_\mathrm{H}=1`$ by $`z7`$, whilst for the DS94 model $`n_\gamma /n_\mathrm{H}=1`$ is achieved at $`z6`$. In the DSGN98 model $`n_\gamma /n_\mathrm{H}=1`$ has not been achieved even by $`z=0`$. When we include recombinations in the IGM, the model with constant gas escape fraction reaches $`F_{\mathrm{fill}}=1`$ only by $`z=6.7`$ for $`f_{\mathrm{clump}}=1`$, and by $`z=5.1`$ for $`f_{\mathrm{clump}}=f_{\mathrm{clump}}^{(\mathrm{halos})}`$, as shown in Fig. 10, showing how reionization is delayed.
Note that while the cosmology considered here is similar to Model G of ?, the parameters of the semi-analytic model used are somewhat different. Specifically, ? used a model in which feedback was much more effective in low mass halos than in our model, since they required their models to produce a B-band luminosity function with a shallow faint end slope. As a result, the epoch at which $`n_\gamma /n_\mathrm{H}=1`$ was much later in Model G of ? than in our current model.
In summary, we see that, even for a specific model of galaxy formation, the predicted epoch of reionization is sensitive to the uncertain values of the escape fraction $`f_{\mathrm{esc}}`$ and the clumping factor $`f_{\mathrm{clump}}`$. If the clumping factor is as large as $`f_{\mathrm{clump}}^{(\mathrm{halos})}`$, then in the case of a constant gas escape fraction $`f_{\mathrm{esc},\mathrm{gas}}`$, we need $`f_{\mathrm{esc},\mathrm{gas}}>10\%`$ in our model to ionize the IGM by $`z=5`$, if absorption by dust is included, and $`f_{\mathrm{esc},\mathrm{gas}}>4\%`$ if dust is ignored. With the more physically-motivated DS94 and DSGN98 models, and the same clumping factor, at most 76% of the IGM is reionized by $`z=5`$, which would be inconsistent with observations of the Gunn-Peterson effect. For the extreme case of a uniform IGM, reionization occurs by $`z=6.1`$ even with the DS94 model for $`f_{\mathrm{esc}}`$. Our “best estimate” is based on combining the DS94 model with $`f_{\mathrm{clump}}^{(\mathrm{halos})}`$ for the IGM clumping factor. As already stated, this model narrowly fails to satisfy the Gunn-Peterson constraint at $`z=5`$ (unless we assume that OB associations are distributed as the gas, rather than lying in the disk mid-plane), suggesting that additional sources of ionizing radiation are required at high redshift, either more stars than in our standard model, or non-stellar sources (e.g. quasars). However, given the theoretical uncertainties in $`f_{\mathrm{esc}}`$ and $`f_{\mathrm{clump}}`$, we consider that this is not yet proven.
## 5 Sensitivity of results to model parameters
We turn now to test the robustness of our results to variations in the parameters of our galaxy formation model. To do this, we have varied key parameters of the models and determined the ionized hydrogen filling factor in each case. We consider several different models. The variant models which we consider are listed in Table 2. In each case, we give the value of the parameter which is changed relative to the standard model given in Table 1.
? have shown that normalising models to the $`z=0`$ B-band luminosity function allows robust estimates of the $`z=0`$ galaxy correlation function to be made. Here we choose a similar constraint, forcing all models to match the $`z=0`$ H$`\alpha `$ luminosity function of ? at $`L_{\mathrm{H}\alpha }=4\times 10^{41}h^2`$ ergs/s (note that at these luminosities the Gallego et al. luminosity function agrees, within the errorbars, with that of ?). This is achieved by adjusting the value of the parameter $`\mathrm{{\rm Y}}`$ (which determines the fraction of brown dwarfs formed in the model). The $`z=0`$ H$`\alpha `$ luminosity functions for all models considered are shown in Fig. 12. Dotted lines show those models with H$`\alpha `$ luminosity functions that are significantly different from that of the standard model (at either the bright or faint ends).
In Table 2 we list escape fractions and filling factors in the variant models for the fixed and DS94 models for $`f_{\mathrm{esc}}`$ for the case $`f_{\mathrm{clump}}=f_{\mathrm{clump}}^{(\mathrm{halos})}`$. Values of $`F_{\mathrm{fill}}`$ may exceed unity, as, in some models, by $`z=5`$ more ionizing photons have escaped into the IGM than are required to reionize the universe.
The standard choice for the feedback efficiency, $`\beta `$, makes feedback highly efficient in galaxies with low circular velocities. In this model $`\beta =f_\mathrm{V}^{\alpha _{\mathrm{hot}}}`$, where $`f_\mathrm{V}=V_{\mathrm{d}isk}/V_{\mathrm{h}ot}`$. The fraction of cold gas which is reheated by supernovae after infinite time (a quantity with direct physical interpretation) is then
$$\frac{\beta }{1R+\beta }=\frac{f_\mathrm{V}^{\alpha _{\mathrm{hot}}}}{1R+f_\mathrm{V}^{\alpha _{\mathrm{hot}}}}.$$
(8)
Thus as $`V_{\mathrm{disk}}0`$ all gas is reheated and no stars are formed. For the modified feedback model, we adapt this form such that even in arbitrarily small potential wells not all the gas is reheated by supernovae. We choose
$$\frac{\beta }{1R+\beta }=\frac{a_\mathrm{r}f_\mathrm{V}^{\alpha _{\mathrm{hot}}}}{1R+f_\mathrm{V}^{\alpha _{\mathrm{hot}}}},$$
(9)
where $`a_\mathrm{r}`$ is an adjustable parameter. Now, as $`V_{\mathrm{disk}}0`$ a fraction $`a_\mathrm{r}`$ of gas is reheated, whilst a fraction $`1a_\mathrm{r}`$ forms stars. The standard feedback model is recovered when $`a_\mathrm{r}=1`$. For the alternative feedback models considered here a value of $`a_\mathrm{r}=0.75`$ is used.
It should be noted that the variation having one of the greatest influences on the predicted filling factors is that of Model 20, where we use the alternative form for feedback given above. This model produces an H$`\alpha `$ luminosity function with a very steep faint end slope, since feedback never becomes highly efficient, even in extremely small dark matter halos. More ionizing photons are produced than in the standard model and higher filling factors are achieved.
Other models which alter the strength of feedback (i.e. Models 3, 4, 5 and 6) also cause large changes in the filling factors. Models with weaker feedback (i.e. Models 3 and 6) result in larger filling factors as they allow more star formation to occur in low mass galaxies (these models again producing a steep slope for the faint end of the H$`\alpha `$ luminosity function). The value of $`\mathrm{\Omega }_\mathrm{b}`$ also has a strong influence on the filling factors as demonstrated by Models 1 and 2. Finally, in Models 21 and 22 we consider two alternative values of $`S_2`$ in the DS94 model. These values span the range of uncertainty for the maximum OB association luminosity in our own Galaxy (?). These models demonstrate that the filling factors predicted by the DS94 model are uncertain by a factor of at least 2 simply because of this uncertainty in the value of $`S_2`$. There is, in fact, further uncertainty introduced as it is not clear if $`S_2`$ represents a real cutoff in the luminosity function of OB associations, or merely a turn-over in that function.
All of the models which significantly alter the predicted filling factors are amongst those marked with a $``$ in Table 2, indicating that such models do not reproduce well the $`z=0`$ H$`\alpha `$ luminosity function, and can therefore be discarded as being unrealistic. With these models removed, our predictions for $`F_{\mathrm{fill}}`$ are reasonably robust. Considering all the realistic models we find that for the fixed gas escape fraction of 10%, $`F_{\mathrm{fill}}`$ at $`z=5`$ is $`1.08_{0.19}^{+0.06}`$ (where the value indicates the filling factor in the standard model and the errors show the range found in the realistic variant models). For the DS94 model we find $`F_{\mathrm{fill}}=0.76_{0.40}^{+0.15}`$ (leaving out the models which vary $`S_2`$ we find $`F_{\mathrm{fill}}=0.76_{0.34}^{+0.07}`$). Our conclusion that with the DS94 escape fractions reionization cannot happen by $`z=5`$ if the clumping factor is as large as $`f_{\mathrm{clump}}^{(\mathrm{halos})}`$ remains valid under all realistic parameter variations considered here. On the other hand, if the clumping factor is closer to the case of a uniform IGM, then reionization by $`z=5`$ is possible in the DS94 model (but not in the DSGN98 model).
So far we have considered a single cosmology, namely $`\mathrm{\Lambda }`$CDM. This choice was motivated by the work of ? and ?, who have shown that the semi-analytic model is able to reproduce many features of the observed galaxy population for this cosmology. However, in order to explore the effects of cosmological parameters on reionization, we have also considered a $`\tau `$CDM cosmology, with $`\mathrm{\Omega }=1`$, in which we use model parameters identical to those of ?. We note that ? were unable to match the galaxy correlation function at $`z=0`$ for this cosmology.
We find that in the $`\tau `$CDM model, our basic results are unchanged, i.e. with physical models for the escape fraction, the IGM is reionized by $`z=5`$ only if it is much less clumped than in our halo clumping model with $`f_{\mathrm{clump}}=f_{\mathrm{clump}}^{(\mathrm{halos})}`$. The escape fractions in this cosmology are actually somewhat higher than in the $`\mathrm{\Lambda }`$CDM cosmology (due to a lower amount of gas and dust in galaxies). At $`z=5`$ the DS94 model predicts a mean escape fraction $`16\%`$, whilst the DSGN98 model predicts $`0.1\%`$. However, the filling factors are significantly lower (for example, in the DS94 model $`F_{\mathrm{fill}}=0.24`$ at $`z=5`$ when $`f_{\mathrm{clump}}^{(\mathrm{halos})}`$ is used as compared to $`F_{\mathrm{fill}}=0.76`$ in $`\mathrm{\Lambda }`$CDM). This reflects the fact that many fewer ionizing photons are produced in this cosmology (due to the fact that a stronger feedback is required in order that the model fits the properties of galaxies at $`z=0`$), and that less of the gas has become collisionally ionized in virialised halos in $`\tau `$CDM than in $`\mathrm{\Lambda }`$CDM. The only factor which works in favour of a higher filling factor in $`\tau `$CDM is that the clumping factor is somewhat lower. However, this is not enough to offset the two effects described above.
## 6 Spatial distribution of ionizing sources and CMB fluctuations
### 6.1 Spatial distribution
We now consider the temperature anisotropies imprinted on the microwave background by the IGM following reionization. These depend on the spatial and velocity correlations of the ionized gas. A fully self-consistent calculation of these correlations on the relevant scales would require very high resolution numerical simulations including both gas dynamics and radiative transfer (e.g. ?). No such numerical simulation is yet available with the necessary combination of volume and resolution to calculate the secondary CMB anisotropies on all angular scales of interest. Therefore in this paper, we calculate the spatial and velocity distribution of the ionized gas in an approximate way, by combining our semi-analytical galaxy formation model with a high resolution N-body simulation of the dark matter.
We have used the same $`\mathrm{\Lambda }`$CDM simulation as ?, described in detail by Jenkins et al. 1998, which has $`\mathrm{\Omega }_0=0.3`$, a cosmological constant $`\mathrm{\Lambda }_0=0.7`$, a Hubble constant of $`h=0.7`$ in units of $`100\mathrm{k}m/s/Mpc`$, and which is normalised to produce the observed abundance of rich clusters at $`z0`$ (?). Using the same semi-analytic model as employed here, ? were able to match the observed galaxy two-point correlation function at $`z=0`$ in this cosmology. The simulation has a box of length 141.3 $`h^1`$ Mpc and contains $`256^3`$ dark matter particles, each of mass of $`1.4\times 10^{10}h^1M_{}`$. We identify halos in this simulation using the friends-of-friends (FOF) algorithm with the standard linking length of 0.2, and then populate them with galaxies according to the semi-analytic model. We consider only groups consisting of 10 particles or more, and so resolve dark halos of mass $`1.4\times 10^{11}h^1M_{}`$ or greater. Sources in halos which are unresolved in the simulations can produce a significant fraction of the total ionizing luminosity, according to the semi-analytic models. To circumvent this problem, we add sources in unresolved halos into the simulation in one of two ways. The first method is to place the sources on randomly chosen dark matter particles which do not belong to any resolved halo. An alternative method is to place these sources completely at random within the simulation volume. This makes the unresolved sources completely unclustered and so is an interesting extreme case. As we will be forced to construct toy models to determine which regions of the simulation are ionized, the exact treatment of these unresolved halos will not be of great importance. The number of unresolved halos added to the simulation volume is determined from the Press-Schechter mass function, multiplied by a correction factor of 0.7 to make it match the low mass end of the N-body mass function in $`\mathrm{\Lambda }`$CDM at $`z=3`$.
In order to calculate the correlations between ionized regions that are needed to determine the temperature anisotropies induced in the CMB, a simulation with at least the volume of this one is required. Unfortunately, with present computing resources, this excludes the possibility of an exact calculation of the shape and size of the ionized regions, which would require much higher resolution, and also the inclusion of gas dynamics and radiative transfer. Therefore we have used five toy models to determine which regions of the simulation are ionized, for a given distribution of ionizing sources. These models cover a range of possibilities which is likely to bracket the true case, and provide an estimate of the present theoretical uncertainties.
For each model, we divide the simulation volume into $`256^3`$ cubic cells, resulting in a cell size of $`0.55h^1`$ Mpc. As the gas distribution is not homogeneous, the volume of gas ionized will depend on the density of gas in the ionized region. We assume that the ionizing luminosity from the galaxies in each halo all originates from the halo centre, and that the total mass $`M`$ of gas ionized by each halo is the same as it would be for an IGM which is uniform on large scales, but with small-scale clumping $`f_{\mathrm{clump}}`$, as given by eqn. (4). We add to this the mass of any collisionally ionized gas in the halo. We then calculate the volume of the ionized region around each halo using $`M=\overline{n}_\mathrm{H}m_\mathrm{H}V`$, where $`\overline{n}_\mathrm{H}`$ is the mean IGM density within the volume V. We use several different toy models to calculate the spatial distribution of ionized gas in the simulations. In all cases, the total mass of hydrogen ionized is assumed to be the same as for a homogeneous distribution with the specified clumping factor.
Model A (Growing front model) Ionize a spherical volume around each halo with a radius equal to the ionization front radius for that halo assuming a large-scale uniform distribution of Hi. Since the Hi in the simulation is *not* uniformly distributed, and also because some spheres will overlap, the ionized volume will not contain the correct total mass of Hi. We therefore scale the radius of each sphere by a constant factor, $`f`$, and repeat the procedure. This process is repeated, with a new value of $`f`$ each time, until the correct total mass of Hi has been ionized.
Model B (High density model) In this model we ignore the positions of halos in the simulation. Instead we simply rank the cells in the simulation volume by their density. We then completely ionize the gas in the densest cell. If this has not ionized enough Hi we ionize the second densest cell. This process is repeated until the correct total mass of Hi has been ionized.
Model C (Low density model) As model B, but we begin by ionizing the least dense cell, and work our way up to cells of greater and greater density. This model mimics that of ?.
Model D (Random spheres model) As Model A but the spheres are placed in the simulation entirely at random rather than on the dark matter halos. By comparing to Model A this model allows us to estimate the importance of the spatial clustering of dark matter halos.
Model E (Boundary model) Ionize a spherical region around each halo with a radius equal to the ionization front radius for that halo. This may ionize too much or not enough Hi depending on the density of gas around each source. We therefore begin adding or removing cells at random from the boundaries of the already ionized regions until the required mass of Hi is ionized.
Fig. 13 shows six slices through the N-body simulation. The top left slice shows the density of all gas (which is assumed to trace the dark matter), whilst the other slices show only the density of ionized gas. Model A shows particularly well the correlated nature of the ionizing sources (due to the fact that galaxies form in the high density regions of the dark matter), as the densest regions of the simulations are the ones which have become most highly ionized.
In Fig. 14 we show the filling factor in the N-body simulation for the fixed $`f_{\mathrm{esc},\mathrm{gas}}`$ model with different values of $`f_{\mathrm{esc},\mathrm{gas}}`$, for the case $`f_{\mathrm{clump}}=f_{\mathrm{clump}}^{(\mathrm{halos})}`$. The filling factors calculated from the simulation are always less (for a given value of $`f_{\mathrm{esc},\mathrm{gas}}`$) than those calculated in §4.1 (see Fig. 10). This is because the simulation contains fewer low mass dark matter halos than predicted by the Press-Schechter theory, hence it contains fewer ionizing sources.
### 6.2 CMB Fluctuations
The reionization of the IGM imprints secondary anisotropies on the CMB through Thomson scattering off free electrons (see, for example, ???). These anisotropies result from the spatially varying ionized fraction and from density and velocity variations in the ionized IGM. The calculation of these secondary effects involves correlation functions of density fluctuations and velocity fields which are easily determined in our models. To predict the form of these fluctuations, we first calculate the two-point correlations between ionized gas over the redshift range 3 to 18, assuming that gas in the IGM traces the dark matter density and velocity. To do this, we use the $`256^3`$ grid of ionization fractions, $`x_\mathrm{e}`$, described in §6.1. We determine in each grid cell the value of
$$\zeta =\left[\frac{x_\mathrm{e}\left(1+\delta \right)}{x_\mathrm{e}(1+\delta )}1\right]v_{\mathrm{los}},$$
(10)
where $`\delta `$ is the dark matter overdensity in the cell, $`v_{\mathrm{los}}`$ is the component of the mean dark matter velocity in the cell along the line of sight to a distant observer, and the averaging of $`x_\mathrm{e}(1+\delta )`$ is over all cells in the simulation volume. The dark matter density and velocity in each cell are estimated by assigning the mass and velocity of each dark matter particle to the grid using a cloud-in-cell algorithm.
We then compute the correlation function
$$\xi _{\zeta \zeta }(r)=\zeta (𝐱)\zeta (𝐱+𝐫)_𝐱.$$
(11)
This correlation function is all that is needed to determine the spectrum of fluctuations imprinted in the CMB by the reionization process. A detailed description of how this spectrum is computed is given in Appendix C.
In Fig. 15 we show the secondary CMB anisotropies calculated as described above. The left-hand panel shows the results for Model E with a fixed gas escape fraction of 0.10, for the cases that unresolved halos are placed either on ungrouped particles (solid line) or at random in the simulation volume (dotted line). The clumping factor is $`f_{\mathrm{clump}}=f_{\mathrm{clump}}^{(\mathrm{halos})}`$, as will be used for all models considered in this section. The particular choice of $`f_{\mathrm{clump}}`$ is not important for our conclusions about the CMB fluctuations, since we will consider different values of $`f_{\mathrm{esc},\mathrm{gas}}`$. The choice of placement scheme is seen to make little difference to the results, the two curves differing by $`<10\%`$ for $`250<\mathrm{}<5000`$, with the difference growing to 40% by $`\mathrm{}=20000`$. The size of the grid cell used in our calculation of the ionized gas correlation function corresponds to $`\mathrm{}25000`$ at $`z=3`$, and the turnover around $`\mathrm{}=10^4`$ is simply due to the cell size. As such, our method is unable to determine the form of the CMB fluctuations at higher $`\mathrm{}`$.
The right-hand panel of Fig. 15 shows the variations in our estimates of $`C_{\mathrm{}}`$ which arise from using the five Models A-E. Here the differences between the curves are larger, with Models B and C differing by a factor of $`2.5`$ at $`\mathrm{}=10^4`$. The amplitude of the curves is affected by the strength of the correlations present in each model (e.g. the “high density” model is the most strongly correlated and has the highest amplitude, whilst the “low density” model has the weakest correlations and hence the lowest amplitude). However, the shapes of the curves are all very similar.
Figure 16 examines the effect on the secondary anisotropies of varying the escape fraction $`f_{\mathrm{esc},\mathrm{gas}}`$. The trend is for increasing amplitude of anisotropy with increasing escape fraction (which results in a higher reionization redshift). If, however, we boost the number of photons produced by increasing $`f_{\mathrm{esc},\mathrm{gas}}`$ above 1 (this of course being unphysical, but a simple way of examining the effects of producing more ionizing photons), little further increase in amplitude is seen.
The form expected for the CMB anisotropies produced by patchy reionization has been calculated for a simple model by ?. In this model, reionization of the universe is assumed to begin at some redshift, $`z_\mathrm{i}`$, and is completed (i.e. the filling factor reaches unity) after a redshift interval $`\delta z`$. Sources are assumed to appear at random positions in space and to each ionize a spherical region of comoving radius $`R`$. Once such an ionized region has appeared it remains forever. In this model, the power $`\mathrm{}^2C_{\mathrm{}}/2\pi `$ is predicted to have the form of white noise at small $`\mathrm{}`$, since the ionized regions are uncorrelated.
We compare the simple model of ? with our own results in Fig. 16. Since in our model reionization has no well-defined starting redshift, and ionized regions span a range of sizes, we simply choose values of $`R`$, $`z_\mathrm{i}`$ and $`\delta z`$ in order to match the two models at the peak in the spectrum (even though the position of this peak in our results is an artifact of our simulation resolution, we simply wish to demonstrate here the difference in small $`\mathrm{}`$ slopes between our model and that of ?). The chosen values of $`R=0.85h^1`$Mpc, $`z_\mathrm{i}=11`$ and $`\delta z=5`$ are all plausible for the ionization history and sizes of ionized regions seen in our model (the mean comoving size of regions ranging from $`1.4h^1`$Mpc at $`z=3`$ to $`0.2h^1`$Mpc at $`z=18`$). Note that ? calculate a somewhat different form for the anisotropy spectrum for this same model, in which the amplitude, $`A`$, is roughly half that found by ?, and the peak in the spectrum occurs at slightly higher $`\mathrm{}`$. Despite these differences both ? and ? agree upon the general form of the spectrum (sharp peak plus white noise at small $`\mathrm{}`$), and this is all we are interested in here.
The $`C_{\mathrm{}}`$ declines much more rapidly as $`\mathrm{}0`$ in the ? model than in ours. Note that Model D, the random sphere model, also shows the same behaviour as our other models, indicating that it is not the correlated positions of the ionizing sources in our model which produce the excess power at small $`\mathrm{}`$. If we force all halos in our model to have equal ionized volumes surrounding them, whilst retaining the same total filling factor, we find that the excess power above the white noise spectrum at small $`\mathrm{}`$ remains, so neither is the excess due to the range of ionizing front radii, $`R`$, present in our model. This excess power can therefore be seen to be due to the correlations in gas density and velocity induced by gravity. In fact, if we repeat our calculations but ignore correlations in the gas density field (i.e. we set $`\delta =0`$ everywhere) we find a CMB spectrum which has a slope for small $`\mathrm{}`$ which is much closer to the ? white-noise slope, and which has an amplitude over five times lower than when density correlations are included. The remaining differences between our model and that of ? in this case are due to the correlated nature of ionized regions in our model.
The amplitude of the secondary anisotropies also depends upon the assumed value of $`\mathrm{\Omega }_\mathrm{b}`$ (as this determines the optical depth for electron scattering). The preferred value for our galaxy formation model of $`\mathrm{\Omega }_\mathrm{b}=0.02`$ is relatively low compared to estimates based on light element abundances and big bang nucleosynthesis. Fig. 17 shows the effect of increasing $`\mathrm{\Omega }_\mathrm{b}`$ to 0.04. If the evolution of the ionized regions were the same in both models, we would expect the amplitude to increase by a factor of four (since it is proportional to the square of the baryon density). The evolution of ionized regions is actually quite similar in the two cases, and so the factor of four increase is seen. With this higher value of $`\mathrm{\Omega }_\mathrm{b}`$, the secondary anisotropies due to patchy reionization would be potentially detectable above $`\mathrm{}3000`$.
We note that a similar approach to computing the spectrum of secondary anisotropies due to patchy reionization has been taken by ?. Using a different model of galaxy formation, ? grow spherical ionization fronts around dark matter halos identified in an N-body simulation, and from these they estimate the spectrum of CMB anisotropies produced. The simulations employed by ? have higher resolution (but much smaller volume) than the GIF simulations used in our work. ? therefore do not have the problem of locating unresolved halos in their simulation, but their calculation of secondary anisotropies is restricted to smaller angular scales (roughly $`5\times 10^3<\mathrm{}<2\times 10^5`$) compared to ours. Furthermore, ? make some approximations in calculating the anisotropies which we do not, ignoring variations in the total IGM density, and assuming that the ionized fraction is completely uncorrelated with the velocity field. ? consider only a single model for reionizing the simulation volume. As we have shown, our five toy models for the distribution of ionized regions lead to factors of 2–3 difference in the secondary anisotropy amplitudes, indicating that the results are not very sensitive to the model adopted for the distribution of ionized regions or to the treatment of unresolved halos in the simulation. ? carried out their calculations in a different cosmology to ours, also with a different value of $`\mathrm{\Omega }_\mathrm{b}`$, but once these differences are taken into account, their results seem reasonably consistent with ours.
## 7 Discussion and Conclusions
We have outlined an approach to studying the reionization of the universe by the radiation from stars in high redshift galaxies. We have focussed on the reionization of hydrogen, but the approach can be generalised to study helium reionization (e.g. ?), and also to include radiation from quasars. Our main conclusions are:
(i) Using a model of galaxy formation constrained by several observations of the local galaxy population, enough ionizing photons are produced to reionize the universe by $`z=11.7`$. This assumes that all ionizing photons escape from the galaxies they originate in, and that the density of the IGM is uniform. Reionization is delayed until $`z10.9`$ in the case of a clumped IGM, in which gas falls into halos with virial temperatures exceeding $`10^4`$K. Galaxies can reionize such a clumped IGM by $`z=5`$ providing that, on average, at least 4% of ionizing photons can escape from the galaxies where they are produced. In the case of a uniform IGM, an escape fraction of only 1.4% is sufficient to reionize by $`z=5`$.
Using a physical model for the escape of ionizing radiation from galaxies, in which photons escape through “Hii chimneys” ionized in the gas layers in galaxy disks (?, hereafter DS94), we predict reionization by $`z=6.1`$ for a uniform IGM or by $`z=4.5`$ for a clumped IGM. Models which assume that all the gas in galaxy disks remains neutral are unable to reionize even a uniform IGM by $`z=0`$. Using alternative estimates of the IGM clumping factor from ? or ?, we find reionization redshifts comparable with those found using our own clumping model, i.e. in the range $`z=`$ 4.5–5.0 with the DS94 model for the escape fraction.
(ii) Once the ionizing escape fraction and IGM clumping factor have been specified, our estimates for the filling factor of ionized gas in the IGM are reasonably robust, providing that we consider only models which are successful in matching the H$`\alpha `$ luminosity function of galaxies at $`z=0`$. By far the greatest remaining influences on the ionized filling factor come from the value of the baryon fraction $`\mathrm{\Omega }_\mathrm{b}`$ and the prescription for feedback from supernovae. However, we have shown that altering these parameters also produces large changes in the $`z=0`$ H$`\alpha `$ luminosity function.
(iii) We combined our model for reionization with N-body simulations of the dark matter distribution in order to predict the spectrum of secondary anisotropies imprinted on the CMB by the process of reionization. The shape of this spectrum is almost independent of the assumptions about reionization, but the amplitude depends on the spatial distribution of the ionized regions, the redshift at which reionization occurs and the baryon fraction. We find considerably more power in the anisotropy spectrum at small $`\mathrm{}`$ than predicted by models which do not account for the large-scale correlations in the gas density and velocity produced by gravity. Despite the uncertainty in the spatial distribution of ionized regions, we are able to determine the amplitude of this spectrum to within a factor of three for a given $`\mathrm{\Omega }_\mathrm{b}`$ (the amplitude being proportional to $`\mathrm{\Omega }_\mathrm{b}^2`$). The results found by ? using a similar technique are reasonably consistent with ours, once differences in $`\mathrm{\Omega }_\mathrm{b}`$ and other cosmological parameters are allowed for.
Detection of these secondary anisotropies, which would constrain the reionization history of the Universe, would require fractional temperature fluctuations of $`10^7`$ to be measured on angular scales smaller than several arcminutes. Although the Planck and MAP space missions are unlikely to have sufficient sensitivity to observe such anisotropies, the Atacama Large Millimeter Array is expected to be able to measure temperature fluctuations of the level predicted at $`\mathrm{}10^4`$ in a ten hour integration.
Previous studies of reionization have either used an approach similar to our own, i.e. employing some type of analytical or semi-analytical model (e.g. ????), or else have used direct hydrodynamical simulations (e.g. ?). While the latter technique can in principle follow the detailed processes of galaxy formation, gas dynamics and radiative transfer, in practice the resolutions attainable at present do not allow such simulations to resolve the small scales relevant to this problem. Furthermore, the implementation of star formation and feedback in such models is far from straightforward.
There are two main uncertainties in our approach, as in most others: the fraction $`f_{\mathrm{esc}}`$ of ionizing photons that escape from galaxies, and the clumping factor $`f_{\mathrm{clump}}`$ of gas in the IGM. Future progress depends on improving estimates of the effects of clumping using larger gas dynamical simulations, on better modelling of the escape of ionizing photons from galaxies, and on better understanding of star formation and supernova feedback in high redshift objects.
## Acknowledgements
AJB and CGL acknowledge receipt of a PPARC Studentship and Visiting Fellowship respectively. AN is supported by a grant from the Israeli Science Foundation. NS is supported by the Sumitomo Foundation and acknowledges the Max Planck Institute for Astrophysics for their warm hospitality. This work was supported in part by a PPARC rolling grant, by a computer equipment grant from Durham University and by the European Community’s TMR Network for Galaxy Formation and Evolution. We acknowledge the Virgo Consortium and GIF for making available the GIF simulations for this study. We are grateful to Martin Haehnelt and Tom Abel for stimulating conversations. We also thank Shaun Cole, Carlton Baugh and Carlos Frenk for allowing us to use their galaxy formation model, and for advice on implementing modifications in that model.
## Appendix A Calculation of the escaping fraction
### A.1 Escaping fraction in the DS94 model: Stars in mid-plane
In the model of ?, hereafter DS94, ionizing photons escape from galactic disks through “Hii chimneys”, which are holes in the neutral gas layer ionized by OB associations. The OB associations are assumed to lie in the disk mid-plane, and to have a distribution of ionizing luminosities $`\mathrm{d}N/\mathrm{d}SS^2`$ for $`S_1<S<S_2`$, $`(\mathrm{d}N/\mathrm{d}S)\mathrm{d}S`$ being the number of associations with luminosities in the range $`S`$ to $`S+\mathrm{d}S`$ (?). The gas is assumed to have a Gaussian vertical distribution with scaleheight $`h_\mathrm{z}`$. The fraction of Lyc photons escaping through chimneys on both sides of the disk at radius $`r`$ is (?, eqn. 24)
$$f_{\mathrm{esc},\mathrm{gas}}=\{\begin{array}{cc}0\hfill & \text{if }S_\mathrm{m}S_2\hfill \\ \left[\mathrm{ln}\left(\frac{S_2}{S_\mathrm{m}}\right)+\frac{9}{2}\left(\frac{S_\mathrm{m}}{S_2}\right)^{1/3}\frac{S_\mathrm{m}}{2S_2}4\right]/\mathrm{ln}\left(\frac{S_2}{S_1}\right)\hfill & \text{if }S_1S_\mathrm{m}<S_2\hfill \\ 1+\left[\frac{9}{2}\left\{\left(\frac{S_\mathrm{m}}{S_2}\right)^{1/3}\left(\frac{S_\mathrm{m}}{S_1}\right)^{1/3}\right\}\frac{1}{2}\left\{\frac{S_\mathrm{m}}{S_2}\frac{S_\mathrm{m}}{S_1}\right\}\right]/\mathrm{ln}\left(\frac{S_2}{S_1}\right)\hfill & \text{if }S_\mathrm{m}<S_1,\hfill \end{array}$$
(12)
where $`S_\mathrm{m}`$ is defined as
$$S_\mathrm{m}(r)=\pi ^{3/2}n_0^2\mathrm{exp}\left(2r/r_{\mathrm{disk}}\right)h_\mathrm{z}^3\alpha _\mathrm{H}^{(2)},$$
(13)
Here $`\alpha _\mathrm{H}^{(2)}`$ is the recombination coefficient for hydrogen for recombinations to all energy levels except the first, and we have assumed an exponential disk with radial scalelength $`r_{\mathrm{disk}}`$, so that the hydrogen gas density is
$$n(r,z)=n_0\mathrm{exp}(r/r_{\mathrm{disk}}z^2/2h_\mathrm{z}^2),$$
(14)
$`r`$ and $`z`$ being the usual cylindrical polar coordinates.
Since $`S_\mathrm{m}`$ varies throughout the galactic disk we average the escape fraction over the entire disk, assuming that the local rate of star formation is proportional to the column density of the disk (??) and that $`h_\mathrm{z}`$ is constant with radius.
The fraction of all ionizing photons produced by the galaxy which can escape into the IGM is then given by,
$$f_{\mathrm{esc},\mathrm{gas}}=\{\begin{array}{cc}\left[\left\{\frac{352}{225}\frac{4}{15}\mathrm{ln}\left(\frac{S_2}{S_\mathrm{m}^0}\right)\right\}\left(\frac{S_2}{S_\mathrm{m}^0}\right)^{1/2}\left\{\frac{352}{225}\frac{4}{15}\mathrm{ln}\left(\frac{S_1}{S_\mathrm{m}^0}\right)\right\}\left(\frac{S_1}{S_\mathrm{m}^0}\right)^{1/2}\right]/\mathrm{ln}\left(\frac{S_2}{S_1}\right)\hfill & \text{if }S_\mathrm{m}^0S_2\hfill \\ \left[\left\{\frac{4}{15}\mathrm{ln}\left(\frac{S_1}{S_\mathrm{m}^0}\right)\frac{352}{225}\right\}\left(\frac{S_1}{S_\mathrm{m}^0}\right)^{1/2}\frac{1}{18}\frac{S_\mathrm{m}^0}{S_2}+\frac{81}{50}\left(\frac{S_\mathrm{m}^0}{S_2}\right)^{1/3}+\mathrm{ln}\left(\frac{S_2}{S_\mathrm{m}^0}\right)\right]/\mathrm{ln}\left(\frac{S_2}{S_1}\right)\hfill & \text{if }S_1S_\mathrm{m}^0<S_2\hfill \\ 1+\left[\frac{81}{50}\left\{\left(\frac{S_\mathrm{m}^0}{S_2}\right)^{1/3}\left(\frac{S_\mathrm{m}^0}{S_1}\right)^{1/3}\right\}\frac{1}{18}\left\{\frac{S_\mathrm{m}^0}{S_2}\frac{S_\mathrm{m}^0}{S_1}\right\}\right]/\mathrm{ln}\left(\frac{S_2}{S_1}\right)\hfill & \text{if }S_\mathrm{m}^0<S_1,\hfill \end{array}$$
(15)
where $`S_\mathrm{m}^0`$ is the value of $`S_\mathrm{m}`$ calculated for $`r=0`$. In Figure 18 we show this average escape fraction as a function of the ratio $`S_\mathrm{m}^0/S_2`$ for $`S_2/S_1=1000`$.
Our model of galaxy formation calculates the radial scale length of each galaxy’s disk and also the mass of cold gas present in that disk, and we assume that $`h_\mathrm{z}/r_{\mathrm{disk}}`$ is constant. We can therefore determine $`n_0`$ and the ratios $`S_\mathrm{m}^0/S_2`$ and $`S_\mathrm{m}^0/S_1`$. Hence $`f_{\mathrm{esc},\mathrm{gas}}`$ can be found using eqn. (15).
### A.2 Escaping fraction in DS94 model: Stars tracing gas
In the DS94 model OB associations are assumed to lie in the midplane of the galaxy disk. If instead OB associations are spread throughout the gas layer, having the same vertical distribution as the cold gas, then the resulting escape fraction will be higher than that in the DS94 model. We assume the same density profile as before, given by eqn.(14). Consider an OB association emitting $`S`$ ionizing photons per second, at position $`(r,z)`$ in the disk. We make the assumption (as did DS94) that the radial variations in density can be ignored for calculating the escape fraction at radius $`r`$ (which will be a valid assumption provided the size of the Hii region formed is much less than $`r_{\mathrm{disk}}`$). In order for any photons emitted into a cone of solid angle $`\mathrm{d}\mathrm{\Omega }`$ which makes an angle $`\theta `$ with the $`z`$-axis to escape the galaxy, the emission rate of photons into this cone must exceed the total recombination rate in the cone. This occurs for an ionizing luminosity $`S_{\mathrm{req}}(\theta )`$, where,
$$S_{\mathrm{req}}(\theta )\frac{\mathrm{d}\mathrm{\Omega }}{4\pi }=n_0^2\alpha _\mathrm{H}^{(2)}\mathrm{exp}\left(2\frac{r}{r_{\mathrm{disk}}}\right)_0^{\mathrm{}}l^2\mathrm{exp}\left(\frac{(z+l\mathrm{cos}\theta )^2}{h_\mathrm{z}^2}\right)dld\mathrm{\Omega },$$
(16)
which can be written as $`S_{\mathrm{req}}^\pm (\theta )=\pm S_{\mathrm{req}}^{0,\pm }/\mathrm{cos}^3\theta `$ ($`S_{\mathrm{req}}^+`$ is the solution for $`\mathrm{cos}\theta >0`$ and $`S_{\mathrm{req}}^{}`$ is the solution for $`\mathrm{cos}\theta <0`$), where
$$S_{\mathrm{req}}^{0,\pm }=S_\mathrm{m}^0\mathrm{e}^{2r/r_{\mathrm{disk}}}\left\{\left[1\mathrm{erf}\left(\frac{z}{h_\mathrm{z}}\right)\right]\left(1+\frac{2z^2}{h_\mathrm{z}^2}\right)\frac{2}{\sqrt{\pi }}\frac{z}{h_\mathrm{z}}\mathrm{exp}\left(\frac{z^2}{h_\mathrm{z}^2}\right)\right\}.$$
(17)
This defines two critical angles, $`\mathrm{cos}\theta _\mathrm{c}^\pm (S)=\pm (S_{\mathrm{req}}^{0,\pm }/S)^{1/3}`$, such that photons can escape the galaxy only if $`\theta <\theta _\mathrm{c}^+(S)`$ or $`\theta >\theta _\mathrm{c}^{}(S)`$.
The total escaping fraction from this OB association is then given by,
$`f_{\mathrm{esc},\mathrm{gas}}(S)`$ $`=`$ $`{\displaystyle \frac{1}{2S}}\left[{\displaystyle _{\mathrm{cos}\theta _\mathrm{c}^+}^1}\left[SS_{\mathrm{req}}^+(\theta )\right]\mathrm{d}(\mathrm{cos}\theta )+{\displaystyle _1^{\mathrm{cos}\theta _\mathrm{c}^{}}}\left[SS_{\mathrm{req}}^{}(\theta )\right]\mathrm{d}(\mathrm{cos}\theta )\right]`$ (18)
$`f_{\mathrm{esc},\mathrm{gas}}(S)`$ $`=`$ $`1{\displaystyle \frac{3}{4}}\left({\displaystyle \frac{S_{\mathrm{req}}^{0,+}}{S}}\right)^{1/3}{\displaystyle \frac{3}{4}}\left({\displaystyle \frac{S_{\mathrm{req}}^{0,}}{S}}\right)^{1/3}+{\displaystyle \frac{S_{\mathrm{req}}^{0,+}}{4S}}+{\displaystyle \frac{S_{\mathrm{req}}^{0,}}{4S}}.`$ (19)
Averaging this escape fraction over the assumed OB association luminosity function then gives a mean escape fraction of
$$f_{\mathrm{esc},\mathrm{gas}}=\{\begin{array}{cc}0\hfill & \text{if }S_2>S_\mathrm{m}^0\mathrm{e}^{2r/r_{\mathrm{disk}}}\hfill \\ \begin{array}{c}\hfill [\frac{1}{2}\mathrm{ln}\left(\frac{S_2^2}{S_{\mathrm{req}}^{0,+}S_{\mathrm{req}}^{0,}}\right)+\frac{9}{4}\{\left(\frac{S_{\mathrm{req}}^{0,+}}{S_2}\right)^{1/3}+\left(\frac{S_{\mathrm{req}}^{0,}}{S_2}\right)^{1/3}2\}\\ \hfill \frac{1}{4}\{\frac{S_{\mathrm{req}}^{0,+}}{S_2}+\frac{S_{\mathrm{req}}^{0,}}{S_2}2\}]/\mathrm{ln}(\frac{S_2}{S_1})\end{array}\hfill & \text{if }S_1<S_\mathrm{m}^0\mathrm{e}^{2r/r_{\mathrm{disk}}}S_2\hfill \\ \begin{array}{c}\hfill 1+[\frac{9}{4}\{\left(\frac{S_{\mathrm{req}}^{0,+}}{S_2}\right)^{\frac{1}{3}}\left(\frac{S_{\mathrm{req}}^{0,+}}{S_1}\right)^{\frac{1}{3}}+\left(\frac{S_{\mathrm{req}}^{0,}}{S_2}\right)^{\frac{1}{3}}\left(\frac{S_{\mathrm{req}}^{0,}}{S_1}\right)^{\frac{1}{3}}\}\\ \hfill \frac{1}{4}\{\frac{S_{\mathrm{req}}^{0,+}}{S_2}\frac{S_{\mathrm{req}}^{0,+}}{S_1}+\frac{S_{\mathrm{req}}^{0,}}{S_2}\frac{S_{\mathrm{req}}^{0,}}{S_1}\}]/\mathrm{ln}(\frac{S_2}{S_1})\end{array}\hfill & \text{if }S_\mathrm{m}^0\mathrm{e}^{2r/r_{\mathrm{disk}}}S_1\hfill \end{array}$$
(20)
This expression is then averaged over the galaxy disk, assuming a star formation rate proportional to the local gas density, to derive the mean escaping fraction for the entire galaxy. This must be done numerically
Although we have ignored radial variations in the density of the gas when computing the escaping fraction from a single OB association, we find by numerical solution that these variations make only a very small difference to the value of $`f_{\mathrm{esc},\mathrm{gas}}`$, and then only for small $`S_\mathrm{m}^0/S_2`$ (see Fig. 18).
### A.3 Escaping fraction in the DSGN98 model
In the model of ?, DSGN98 the stars producing the Lyman continuum photons are assumed to be uniformly mixed with the gas in the galaxy, which is distributed in an exponential disk. All of the hydrogen in the galaxy is assumed to be in the form of Hi, allowing the optical depth for ionizing photons to be calculated.
We have calculated the escaping fraction in this model exactly, using the density profile given by eqn. (14). As in the case of the DS94 model with stars mixed uniformly with the gas, we begin by finding the escaping fraction as a function of position $`(r_0,z_0)`$ and line of sight $`(\theta ,\varphi )`$. For the DSGN98 model we therefore find the total optical depth in neutral hydrogen along the line of sight, which is
$$\tau _(r_0,z_0,\theta ,\varphi )=\sigma _{\mathrm{H}\mathrm{i}}^{\mathrm{}}_0n(r,z)\mathrm{d}l,$$
(21)
where $`\sigma _{\mathrm{H}\mathrm{i}}=6.3\times 10^{18}\mathrm{cm}^2`$ is the cross-section for hydrogen ionization at 912Å. The quantity $`\mathrm{exp}(\tau )`$ is then averaged over all $`r_0`$, $`z_0`$, $`\theta `$ and $`\varphi `$, assuming that the star formation rate is proportional to the local gas density, to obtain the final escaping fraction. The average escape fraction for the entire galaxy depends only on the quantity $`\tau _0=n_0r_{\mathrm{disk}}\sigma _\mathrm{H}`$ for a given value of $`h_\mathrm{z}/r_{\mathrm{disk}}`$, and is shown in Fig. 19 for $`h_\mathrm{z}/r_{\mathrm{disk}}=0.1`$.
### A.4 Escaping fraction in starbursts
In the case of a burst of star formation triggered by a major merger, we use the same $`f_{\mathrm{esc},\mathrm{gas}}`$ as for quiescent star formation in the the case where $`f_{\mathrm{esc},\mathrm{gas}}`$ is assumed fixed, but in the DS94 and DSGN98 models we estimate the escape fraction by assuming the burst has an approximately spherical geometry, throughout which star formation proceeds uniformly. We assume a sphere of uniform hydrogen number density, $`n`$, given by
$$n=\frac{3M_{\mathrm{gas}}}{4\pi r_{\mathrm{burst}}^31.4\mathrm{m}_\mathrm{H}},$$
(22)
where $`M_{\mathrm{gas}}`$ is the mass of cold gas in the burst, $`r_{\mathrm{burst}}`$ is the radius of the region in which the burst occurs, and $`\mathrm{m}_\mathrm{H}`$ is the mass of a hydrogen atom. The factor of 1.4 accounts for the presence of helium in the gas. We assume also that photons escape only from an outer shell of thickness $`l`$, within which the optical depth is less than 1. Therefore,
$$nl\sigma _{\mathrm{H}\mathrm{i}}1,$$
(23)
where $`\sigma _{\mathrm{H}\mathrm{i}}`$ is the cross section for hydrogen ionization. The escape fraction is simply the fraction of the sphere’s volume in this shell, i.e.
$$f_{\mathrm{esc},\mathrm{gas}}\frac{4\pi r_{\mathrm{burst}}^2l}{4\pi /3r_{\mathrm{burst}}^3}.$$
(24)
Substituting for $`l`$ then gives
$`f_{\mathrm{esc},\mathrm{gas}}`$ $`=`$ $`{\displaystyle \frac{3}{r_{\mathrm{burst}}n\sigma _{\mathrm{H}\mathrm{i}}}}`$ (25)
$`=`$ $`{\displaystyle \frac{4\pi r_{\mathrm{burst}}^21.4\mathrm{m}_\mathrm{H}}{M_{\mathrm{gas}}\sigma _{\mathrm{H}\mathrm{i}}}}.`$
We take $`r_{\mathrm{burst}}`$ to be equal to $`0.1r_{\mathrm{bulge}}`$ where $`r_{\mathrm{bulge}}`$ is the half-mass radius of the bulge formed by the merger. This choice is motivated by observational fact which shows that starburst activity is usually confined to the nuclear region, the size of which is much smaller than that of the galaxy as a whole (e.g. ? and references therein). ? have carried out more elaborate calculations of escaping fractions from spherical galaxies. However, their results are applicable to gas in hydrostatic equilibrium with an NFW dark matter profile and so are not well suited to the case of starbursts.
The star formation rate in the burst is assumed to decline exponentially, with an e-folding time equal to $`f_{\mathrm{dyn}}`$ times the bulge dynamical time. Unless noted otherwise, we assume $`f_{\mathrm{dyn}}=1`$ in all models. As the burst proceeds the mass of cold gas present, $`M_{\mathrm{gas}}`$, declines as it is turned into stars. The escape fraction, $`f_{\mathrm{esc},\mathrm{gas}}`$, therefore increases during the burst, reaching unity as the amount of gas present drops to zero. However, as the star formation rate is declining exponentially during the burst only a small fraction of photons are produced whilst $`f_{\mathrm{esc},\mathrm{gas}}`$ is high.
## Appendix B Calculation of Clumping Factor
To estimate the clumping factor of the photoionized IGM, we make the simplifying assumption that gas in the universe can be split into three components — that which has fallen into dark matter halos and is collisionally ionized or is part of a galaxy, that which has fallen into dark matter halos and is *not* collisionally ionized, and that which has remained outside halos and is smoothly distributed. The first component makes no contribution to the clumping factor. We define the clumping factor as
$$f_{\mathrm{clump}}=\frac{\rho _{\mathrm{IGM}}^2}{\overline{\rho }_{\mathrm{IGM}}^2}=\frac{\rho _{\mathrm{IGM}}^2}{f_{\mathrm{IGM}}^2\overline{\rho }^2},$$
(26)
where $`\rho _{\mathrm{IGM}}`$ is the IGM gas density at any point in the universe (i.e. it does *not* include contributions from collisionally ionized gas or galaxies), $`\overline{\rho }_{\mathrm{IGM}}=f_{\mathrm{IGM}}\overline{\rho }`$ is the mean density of gas in the IGM and $`\overline{\rho }`$ is the mean density of all gas in the universe (here $`f_{\mathrm{IGM}}`$ is the fraction of the total mass of gas in the universe which resides in the IGM, as defined in §2.2).
Let $`f_{\mathrm{m},\mathrm{clumped}}`$ be the fraction of mass in halos above the Jeans halo mass, $`M_\mathrm{J}`$, as calculated from the Press-Schechter mass function for example. These halos occupy a fraction of the volume of the universe given by $`f_{\mathrm{v},\mathrm{clumped}}=f_{\mathrm{m},\mathrm{clumped}}/\mathrm{\Delta }_{\mathrm{vir}}`$. Here, $`\mathrm{\Delta }_{\mathrm{vir}}`$ is the mean density within the virial radius of a halo in units of the mean density of the Universe. The smooth component of gas is assumed to uniformly fill the region outside halos with $`M>M_\mathrm{J}`$, and so has density
$$\rho _{\mathrm{smooth}}=\overline{\rho }f_{\mathrm{m},\mathrm{smooth}}/f_{\mathrm{v},\mathrm{smooth}},$$
(27)
where $`f_{\mathrm{m},\mathrm{smooth}}=1f_{\mathrm{m},\mathrm{clumped}}`$ is the mass fraction of gas in this smooth component, and $`f_{\mathrm{v},\mathrm{smooth}}=1f_{\mathrm{v},\mathrm{clumped}}`$ is the fraction of the volume of the universe that it occupies.
Consider next the non-collisionally ionized gas in a *single* dark matter halo. Averaging over the volume of this one halo we obtain
$$\rho _{\mathrm{clumped}}^2=f_{\mathrm{int}}(1f_{\mathrm{gal}})^2(1x_\mathrm{H})^2\mathrm{\Delta }_{\mathrm{vir}}^2\overline{\rho }^2,$$
(28)
where $`f_{\mathrm{gal}}`$ is the fraction of the baryons which have become part of galaxies within the halo, $`x_\mathrm{H}`$ is the ionized fraction for the hydrogen in the halo gas assuming collisional ionization equilibrium (which we take from the calculations of ?), and $`f_{\mathrm{int}}`$ is a factor of order unity which depends on the shape of the halo gas density profile and is given by
$$f_{\mathrm{int}}=\frac{_0^{r_{\mathrm{vir}}}\rho ^2(r)r^2dr}{_0^{r_{\mathrm{vir}}}\overline{\rho }_{int}^2r^2dr}.$$
(29)
Here $`r_{\mathrm{vir}}`$ is the virial radius of the halo, $`\rho (r)`$ is the density profile of the diffuse gas in the halo, and $`\overline{\rho }_{int}`$ is the mean density of this gas within the virial radius. We ignore any dependence of the density profile of the gas in the halo on the fraction which has cooled to form galaxies. Our results should be insensitive to this assumption, as $`f_{\mathrm{gal}}1`$ in halos where $`x_\mathrm{H}`$ is significantly less than unity.
To find the contribution of gas in halos to the clumping factor, we integrate the above expression over all halos more massive than $`M_\mathrm{J}`$, weighting by the volume for each halo. Adding the contribution from the smooth component, we then obtain
$$f_{\mathrm{clump}}=\frac{f_{\mathrm{m},\mathrm{smooth}}^2}{f_{\mathrm{v},\mathrm{smooth}}f_{\mathrm{IGM}}^2}+\frac{f_{\mathrm{int}}\mathrm{\Delta }_{\mathrm{vir}}}{f_{\mathrm{IGM}}^2}_{M_\mathrm{J}}^{\mathrm{}}(1f_{\mathrm{gal}})^2(1x_\mathrm{H})^2\frac{M_{\mathrm{halo}}}{\rho _\mathrm{c}\mathrm{\Omega }_0}\frac{\mathrm{d}n}{\mathrm{d}M_{\mathrm{halo}}}dM_{\mathrm{halo}},$$
(30)
where we have used that fact that the comoving volume of a dark matter halo of mass $`M_{\mathrm{halo}}`$ is $`M_{\mathrm{halo}}/(\mathrm{\Delta }_{\mathrm{vir}}\mathrm{\Omega }_0\rho _\mathrm{c})`$ ($`\rho _\mathrm{c}`$ being the critical density of the universe at $`z=0`$). Here $`(1f_{\mathrm{gal}})^2`$ is averaged over all halos of mass $`M_{\mathrm{halo}}`$ in our model of galaxy formation.
We determine $`M_\mathrm{J}`$ by finding the mass of a dark matter halo which has a potential well deep enough that it can just hold onto reionized gas. This gives us the minimum mass halo within which gas collects. For the halo to just retain its gas,
$$\frac{\mathrm{d}P}{\mathrm{d}r}=\frac{\mathrm{G}M_\mathrm{J}}{r_{\mathrm{vir}}^2}\rho (r_{\mathrm{vir}}),$$
(31)
where $`r_{\mathrm{vir}}`$ is the virial radius of the halo and $`P`$ is the gas pressure. We approximate this as
$$\frac{P}{r_{\mathrm{vir}}}\frac{\mathrm{G}M_\mathrm{J}}{r_{\mathrm{vir}}^2}\rho (r_{\mathrm{vir}}),$$
(32)
and using the ideal gas law this becomes
$$\frac{\mathrm{k}_\mathrm{B}T}{\mu \mathrm{m}_\mathrm{H}}\frac{\mathrm{G}M_\mathrm{J}}{r_{\mathrm{vir}}}=\frac{4\pi }{3}\mathrm{G}r_{\mathrm{vir}}^2\rho _\mathrm{c}\mathrm{\Omega }_0\mathrm{\Delta }_{\mathrm{vir}}(1+z)^3,$$
(33)
where we have used the relation $`M_\mathrm{J}=4\pi \rho _\mathrm{c}\mathrm{\Omega }_0(1+z)^3\mathrm{\Delta }_{\mathrm{vir}}r_{\mathrm{vir}}^3/3`$. The virial radius is therefore
$$r_{\mathrm{vir}}=\left(\frac{3}{4\pi }\frac{\mathrm{k}_\mathrm{B}T}{\mathrm{G}\mu \mathrm{m}_\mathrm{H}\rho _\mathrm{c}\mathrm{\Omega }_0\mathrm{\Delta }_{\mathrm{vir}}}\right)^{1/2}(1+z)^{3/2},$$
(34)
and the minimum halo mass in which gas is retained is
$$M_\mathrm{J}=\frac{4\pi }{3}\rho _\mathrm{c}(1+z)^3\mathrm{\Omega }_0\mathrm{\Delta }_{\mathrm{vir}}r_{\mathrm{vir}}^3.$$
(35)
We evaluate $`f_{\mathrm{int}}`$ for the case of an isothermal profile with core radius $`r_\mathrm{c}`$:
$$\rho (r)\frac{1}{r^2+r_\mathrm{c}^2}.$$
(36)
The simulations of galaxy clusters by ? and ? show that the gas density profile is well described by this form. Substituting this in eqn. (29), we find
$$f_{\mathrm{int}}=\frac{1}{6}\left(\frac{r_{\mathrm{vir}}}{r_\mathrm{c}}\right)^3\left[\frac{r_{\mathrm{vir}}}{r_\mathrm{c}}\mathrm{arctan}\frac{r_{\mathrm{vir}}}{r_\mathrm{c}}\right]^2\left[\mathrm{arctan}\frac{r_{\mathrm{vir}}}{r_\mathrm{c}}\frac{r_{\mathrm{vir}}}{r_\mathrm{c}}\left(1+\frac{r_{\mathrm{vir}}^2}{r_\mathrm{c}^2}\right)^1\right].$$
(37)
For a typical value of $`r_{\mathrm{vir}}/r_\mathrm{c}=10`$, we therefore find $`f_{\mathrm{int}}=3.14`$.
## Appendix C The Spectrum of CMB Secondary Anisotropies
In this paper, we concentrate on the kinematic Sunyaev-Zel’dovich effect which is induced by the peculiar motions (deviations from pure Hubble flow) of free electrons in ionized regions (??). There exist other secondary sources of CMB anisotropies. However, on angular scales smaller than a few arc-minutes, the kinematic Sunyaev-Zel’dovich effect is likely to provide a dominant contribution. For example, it is known that the temperature anisotropies caused by non-linear growth of density perturbations, which are often referred to as the Rees-Sciama effect or integrated Sachs-Wolfe effect, are of order $`10^7`$ or less (?). These anisotropies depend on the the baryon bulk physical peculiar velocity, $`𝐯`$, and the number density of free electrons, $`n_\mathrm{e}`$. In our calculations of the anisotropies we assume that the $`𝐯`$ is equal to the bulk velocity of the dark matter and that $`n_\mathrm{e}`$ in ionized regions is proportional to the dark matter density.
The temperature anisotropy $`\mathrm{\Theta }(𝜸)=\frac{\mathrm{\Delta }T}{T}`$ observed in a given line of sight direction $`𝜸`$ is (e.g. ?)
$$\mathrm{\Theta }(𝜸,\eta _0)=_{\eta _{\mathrm{rec}}}^{\eta _0}\frac{\mathrm{d}\eta }{(1+z)}\gamma _iv_\mathrm{B}^i\dot{\tau },$$
(38)
where $`\eta (1+z)dt`$ is conformal time with its values at recombination and present denoted, respectively, by $`\eta _{\mathrm{rec}}`$ and $`\eta _0`$. In eqn. (38) we have assumed an optically thin universe. In an optically thick universe these temperature fluctuations are damped by a factor $`\mathrm{e}^\tau `$, where the optical depth is $`\tau =d\eta \sigma _\mathrm{T}n_\mathrm{e}/(1+z)`$, where $`\sigma _\mathrm{T}`$ is the cross section for Thomson scattering. If the universe became instantaneously fully ionized after some redshift $`z_\mathrm{i}`$, the relation between the optical depth $`\tau (\eta _\mathrm{i},\eta _0)`$ and $`z_\mathrm{i}`$ is approximately obtained as $`z_\mathrm{i}=100\mathrm{\Omega }_0\left(0.025/\mathrm{\Omega }_\mathrm{b}h\right)^{2/3}\tau ^{1/3}`$. Therefore, if the reionization takes place at $`z100\mathrm{\Omega }_0\left(0.025/\mathrm{\Omega }_\mathrm{b}h\right)^{2/3}`$, as is the case in our reionization model, then the damping factor can be neglected.
The usual procedure to obtain the angular correlation function of temperature anisotropies in eqn. (38) is by means of Limber’s equation in Fourier space (see for example ?). However, in this paper, we work in real space since we have the two point correlation functions of density and velocity fields directly measured in real space from N-body simulations.
The temperature angular correlation $`C(\theta )`$ can be written as
$$C(\theta )=\sigma _\mathrm{T}^2_{\eta _{\mathrm{r}ec}}^{\eta _0}d\eta _{\eta _{\mathrm{rec}}}^{\eta _0}d\eta ^{}\gamma _i\gamma _j^{}<v^i(𝐱,\eta )v^j(𝐱^{},\eta ^{})n_\mathrm{e}(𝐱,\eta )n_\mathrm{e}(𝐱^{},\eta ^{})>,$$
(39)
where $`\gamma _i\gamma ^i=\mathrm{cos}\theta `$, and, $`𝐱`$ and $`𝐱^{}`$ refer, respectively, to comoving coordinates in the past light geodesics in the directions $`𝜸`$ and $`𝜸^{\mathbf{}}`$ at $`\eta `$ and $`\eta ^{}`$. We write $`n_\mathrm{e}`$ terms of density fluctuations $`\delta `$ as
$$n_\mathrm{e}(𝐱,\eta )=\overline{n}_\mathrm{e}(\eta )x_\mathrm{e}(𝐱,\eta )\left[1+\delta (𝐱,\eta )\right],$$
(40)
where $`\overline{n}_\mathrm{e}(\eta )`$ is the mean total (free and bound) electron number density at time $`\eta `$, and $`x_\mathrm{e}(𝐱)`$, the ionization fraction, is unity in ionized regions and zero otherwise. The correlation lengths of velocity and density fields are small compared to the Hubble radius so that we can approximate $`n_\mathrm{e}(𝐱^{},\eta ^{})=n_\mathrm{e}(𝐱^{},\eta )`$ and similarly for the $`v`$, in eqn. (39).
$$\zeta =\left[\frac{x_\mathrm{e}\left(1+\delta \right)}{x_\mathrm{e}(1+\delta )}1\right]v_{\mathrm{los}},$$
(41)
where $`v_{\mathrm{los}}=\gamma _iv^i`$ is the velocity component in the direction $`\gamma `$. Therefore $`C(\theta )`$ can be written in terms of the velocity correlation function $`\xi _{vv}(y)<v_{\mathrm{los}}(𝐱)v_{\mathrm{los}}(𝐱+𝐲)>`$ and the density-velocity correlation function $`\xi _{\zeta \zeta }(y)<\zeta (𝐱)\zeta (𝐱+𝐲)>`$, both evaluated for fields at the same $`\eta `$.
$$C(\theta )=\sigma _\mathrm{T}^2_{\eta _{\mathrm{rec}}}^{\eta _0}\frac{\mathrm{d}\eta }{1+z}_{\eta _{\mathrm{rec}}}^{\eta _0}\frac{\mathrm{d}\eta ^{}}{1+z^{}}\overline{n}_\mathrm{e}(\eta )\overline{n}_\mathrm{e}(\eta ^{})x_\mathrm{e}(1+\delta )^2\left[\xi _{\zeta \zeta }(|𝐱^{}𝐱|)+\xi _{vv}(|𝐱^{}𝐱|)+\zeta (𝐱)v_{\mathrm{los}}(𝐱^{})+\zeta (𝐱^{})v_{\mathrm{los}}(𝐱)\right].$$
(42)
The dominant contribution to $`C(\theta )`$ is from the term involving $`\xi _{\zeta \zeta }`$. The integration over $`\xi _{vv}`$ yields to phase cancellation (???). The last term in the integrand also has negligible contribution<sup>1</sup><sup>1</sup>1It is interesting that the contribution from the integral over $`\xi _{\zeta \zeta }`$ is still dominant even if we approximate $`\xi _{\zeta \zeta }=\xi _{\delta \delta }\xi _{vv}`$, i.e. if we ignore any correlations between the density and velocity fields. We have checked that the dominant term produces at least an order of magnitude larger anisotropies than the other terms.
In flat space we use the triangle relation, $`|𝐱^{}𝐱|^2=x^2+x^22xx^{}\mathrm{cos}\theta ,`$ we first carry out the integration of eqn. (42) in terms of $`\eta ^{}`$ for fixed $`\eta `$ and $`\theta `$. We compute $`\xi _{\zeta \zeta }`$ at an average redshift $`\overline{z}`$ given by $`1/(1+\overline{z})=\left(1/(1+z_1)+1/(1+z_2)\right)/2`$, which is an appropriate approximation if the correlation length is negligible relative to the horizon scale. It is straightforward to extend the calculation to an open geometry.
From the temperature angular correlation $`C(\theta )`$, we can obtain $`C_{\mathrm{}}`$ as
$$C_{\mathrm{}}=2\pi _1^1\mathrm{d}\mathrm{cos}\theta P_{\mathrm{}}(\mathrm{cos}\theta )C(\theta ),$$
(43)
where $`P_{\mathrm{}}(\mathrm{cos}\theta )`$ is the Legendre polynomial. |
warning/0002/hep-th0002250.html | ar5iv | text | # Thermodynamics of Field Theories from Spinning Branes Talk presented by N.O. at Quantum aspects of gauge theories, supersymmetry and unification, Paris, France (Sept. 1-7, 1999). NBI-HE-00-11, NORDITA-2000/15 HE, hep-th/0002250
## 1 Introduction
The discovery that four-dimensional black holes have thermodynamic properties due to Hawking radiation has been one of the main motivations to study the thermodynamics of black $`p`$-branes in string and M-theory . Through the recently conjectured correspondence between the near-horizon limit of these brane solutions and quantum field theories in the large $`N`$ limit , it has become clear that this study not only probes the nature of quantum gravity, but also provides information about the thermodynamics of quantum field theories in the large $`N`$ limit . In particular, for the non-dilatonic branes (D3,M2,M5) the near-horizon limit of the supergravity solutions has been conjectured to be dual to a certain limit of the corresponding conformal field theories. For the more general dilatonic branes of type II string theory preserving 16 supersymmetries, similar duality relations have also been obtained . See Ref. for a comprehensive review on the AdS/CFT correspondence and for lectures on branes and AdS space.
A common feature of these dualities between near-horizon backgrounds and field theories, is that the supergravity black $`p`$-brane solution exhibits an $`SO(d)`$ isometry (where $`d=Dp1`$ is the dimension of the transverse space) which manifests itself as the R-symmetry of the dual field theory. As a consequence, by considering spinning black $`p`$-brane solutions, we expect on the one hand to learn more about the field theory side, and on the other hand, be able to perform further non-trivial tests of the duality conjectures that include the dependence on this R-symmetry group. The thermodynamics on the two dual sides, including a stability analysis, turns out to provide a useful starting point for this comparison. For the non-dilatonic branes this study has been initiated in a number of recent papers and was further generalized in for all branes that are half-BPS solutions of string and M-theory. This work was extended in , which includes the first construction of spinning brane bound states along with their thermodynamics, which are relevant for non-commutative field theories. In this talk, we will summarize the main results of Refs. , focusing mainly on the spinning $`p`$-branes, their thermodynamics and application to the supergravity/field theory correspondence, while commenting at the end on the spinning brane bound states.
## 2 Spinning $`p`$-branes and the near-horizon limit
The most general black $`p`$-brane solution can be derived by oxidization from the rotating black hole solutions in . We focus on general spinning black $`p`$-branes that are half-BPS solutions of string and M-theory (in the extremal and non-rotating limit), which include the M2 and M5-branes of M-theory and the D and NS-branes of string theory. The asymptotically-flat solution of a spinning black $`p`$-brane in $`D`$ dimensions takes the form
$`ds^2`$ $`=`$ $`H^{\frac{d2}{D2}}\left(fdt^2+{\displaystyle \underset{i=1}{\overset{p}{}}}(dy^i)^2\right)`$ (1)
$`+H^{\frac{p+1}{D2}}\left(\overline{f}^1K_ddr^2+\mathrm{\Lambda }_{\alpha \beta }d\eta ^\alpha d\eta ^\beta \right)`$
$`+H^{\frac{d2}{D2}}{\displaystyle \frac{1}{K_dL_d}}{\displaystyle \frac{r_0^{d2}}{r^{d2}}}({\displaystyle \underset{i,j=1}{\overset{n}{}}}l_il_j\mu _i^2\mu _j^2d\varphi _id\varphi _j`$
$`2\mathrm{cosh}\alpha {\displaystyle \underset{i=1}{\overset{n}{}}}l_i\mu _i^2dtd\varphi _i)`$
where the harmonic function $`H`$ is given by
$$H=1+\frac{1}{K_dL_d}\frac{r_0^{d2}\mathrm{sinh}^2\alpha }{r^{d2}}$$
(2)
and we refer to for further details, including the corresponding expressions of the dilaton $`\varphi `$ and gauge potential $`A_{p+1}`$. Since the transverse space is $`d=Dp1`$ dimensional, these spinning solutions are characterized by a set of angular momenta $`l_i`$, $`i=1\mathrm{}n`$ where $`n=\mathrm{rank}(SO(d))`$, along with the non-extremality parameter $`r_0`$ and another parameter $`\alpha `$ related to the charge.
Using standard methods of black hole thermodynamics we can compute the relevant thermodynamic quantities of the general solution. The mass $`M`$, charge $`Q`$ and angular momenta $`J_i`$ are determined by the background in the asymptotic region. The Hawking-Bekenstein entropy $`S`$ as well as all the intensive quantities, temperature $`T`$, chemical potential $`\mu `$ and angular velocities $`\mathrm{\Omega }_i`$ are computed from the properties at the horizon. The final results are summarized by
$`M=\rho _dr_0^{d2}\left(d1+(d2)\mathrm{sinh}^2\alpha \right)`$ (3)
$`T`$ $`=`$ $`{\displaystyle \frac{d22\kappa }{4\pi r_H\mathrm{cosh}\alpha }},S=4\pi \rho _dr_0^{d2}r_H\mathrm{cosh}\alpha `$
$`\mu `$ $`=`$ $`\mathrm{tanh}\alpha ,Q=\rho _dr_0^{d2}(d2)\mathrm{sinh}\alpha \mathrm{cosh}\alpha `$
$`\mathrm{\Omega }_i`$ $`=`$ $`{\displaystyle \frac{l_i}{(l_i^2+r_H^2)\mathrm{cosh}\alpha }},J_i=2\rho _dr_0^{d2}l_i\mathrm{cosh}\alpha `$
where $`r_H`$ is the horizon radius determined by
$$\underset{i=1}{\overset{n}{}}\left(1+\frac{l_i^2}{r_H^2}\right)r_H^{d2}=r_0^{d2}$$
(4)
and $`\kappa =_{i=1}^nl_i^2/(l_i^2+r_H^2)`$. We have also defined $`\rho _d=V_pV(S^{d1})/16\pi G`$ with $`V_p`$ the worldvolume of the $`p`$-brane and $`V(S^{d1})`$ the volume of the (unit radius) transverse $`(d1)`$-sphere. The quantities (3) obey the conventional Smarr formula
$$M=\frac{d1}{d2}(TS+\mathrm{\Omega }J)+\mu Q$$
(5)
which follows from the first law of thermodynamics and the scaling behavior of the extensive quantities.
Our main interest, however, is in the near-horizon limit of these spinning branes, defined as follows: One introduces a dimensionfull parameter $`\mathrm{}`$ and performs the rescaling
$`r={\displaystyle \frac{r_{\mathrm{old}}}{\mathrm{}^2}},r_0={\displaystyle \frac{(r_0)_{\mathrm{old}}}{\mathrm{}^2}},l_i={\displaystyle \frac{(l_i)_{\mathrm{old}}}{\mathrm{}^2}}`$
$`h^{d2}={\displaystyle \frac{h_{\mathrm{old}}^{d2}}{\mathrm{}^{2d8}}},G={\displaystyle \frac{G_{\mathrm{old}}}{\mathrm{}^{2(d2)}}}`$ (6)
where the new quantities on the left hand side are expressed in terms of the old quantities labelled with a subscript “old”, and $`h`$ is defined by $`h^{d2}=r_0^{d2}\mathrm{cosh}\alpha \mathrm{sinh}\alpha `$. The near-horizon limit is then defined as the limit $`\mathrm{}0`$ keeping all the new quantities in (2) fixed. The rescaling (2) is accompanied by appropriate rescalings of the fields $`g,\varphi ,A`$ such that the new ones are finite, which turns out to precisely leave the action invariant. The resulting near-horizon spinning $`p`$-brane is
$`ds^2`$ $`=`$ $`H^{\frac{d2}{D2}}\left(fdt^2+{\displaystyle \underset{i=1}{\overset{p}{}}}(dy^i)^2\right)`$
$`+H^{\frac{p+1}{D2}}\left(\overline{f}^1K_ddr^2+\mathrm{\Lambda }_{\alpha \beta }d\eta ^\alpha d\eta ^\beta \right)`$
$`2H^{\frac{d2}{D2}}{\displaystyle \frac{1}{K_dL_d}}{\displaystyle \frac{h^{\frac{d2}{2}}r_0^{\frac{d2}{2}}}{r^{d2}}}{\displaystyle \underset{i=1}{\overset{n}{}}}l_i\mu _i^2dtd\varphi _i`$
where the harmonic function is now given by
$$H=\frac{1}{K_dL_d}\frac{h^{d2}}{r^{d2}}$$
(8)
and we refer to for the corresponding expressions of the dilaton and gauge potential.
The thermodynamics in the near-horizon limit is then
$`E=\rho _d{\displaystyle \frac{d}{2}}r_0^{d2}`$ (9)
$`T={\displaystyle \frac{d22\kappa }{4\pi r_H}}{\displaystyle \frac{r_0^{\frac{d2}{2}}}{h^{\frac{d2}{2}}}},S=4\pi \rho _dr_0^{\frac{d2}{2}}h^{\frac{d2}{2}}r_H`$
$`\mathrm{\Omega }_i={\displaystyle \frac{l_i}{(l_i^2+r_H^2)}}{\displaystyle \frac{r_0^{\frac{d2}{2}}}{h^{\frac{d2}{2}}}},J_i=2\rho _dr_0^{\frac{d2}{2}}h^{\frac{d2}{2}}l_i`$
where the internal energy is obtained from the energy above extremality $`E=MQ`$. Note that since the charge and chemical potential have become constant these are not thermodynamic parameters anymore, so that the thermodynamic quantities are given in terms of the $`n+1`$ supergravity parameters $`(r_0,l_i)`$. Due to the different scaling behavior the near-horizon quantities (9) do not satisfy the Smarr law (5), but instead the near-horizon Smarr law
$$E=\frac{d}{2(d2)}(TS+\mathrm{\Omega }J)$$
(10)
It is not difficult to obtain the energy function of the microcanonical ensemble in terms of the extensive variables using the horizon equation (4) and (9), yielding
$`E^{d/2}`$ $`=`$ $`\left({\displaystyle \frac{d}{2}}\right)^{d/2}\rho _d^{(d4)/2}h^{(d2)^2/2}\left({\displaystyle \frac{S}{4\pi }}\right)^{d2}`$ (11)
$`{\displaystyle \underset{i}{}}(1+\left({\displaystyle \frac{2\pi J_i}{S}}\right)^2)`$
Moreover, for any near-horizon spinning $`p`$-brane solution with $`d`$ transverse dimensions, the Gibbs free energy takes the simple form
$$F=ETS\mathrm{\Omega }J=\rho _d\frac{d4}{2}r_0^{d2}$$
(12)
Except for the case of one non-zero angular momentum, it is not possible in general to obtain a closed-form expression of $`F`$ in terms of the proper intensive quantities, the temperature $`T`$ and the angular velocities $`\mathrm{\Omega }_i`$. However, in a low angular momentum expansion this change of variables can be achieved for the general spinning case to any desired order in $`\omega _i=\mathrm{\Omega }_i/T`$, and will be sufficient for some of our applications below. We also remark that the Gibbs free energy (12) is properly reproduced for all near-horizon spinning branes by computing the on-shell Euclidean action, taking into account the boundary term. This is not only an important consistency check but also essential when calculating string corrections to the free energy of spinning branes.
## 3 Thermodynamics of dual field theories
The thermodynamics of spinning branes and their near-horizon limit is interesting in its own right, but also enables us to obtain information on the thermodynamics of field theories in the strongly coupled large $`N`$ limit, via the near-horizon supegravity/field theory correspondence . The fact that the branes are spinning, introduces the new thermodynamic parameters $`\mathrm{\Omega }_i`$ and $`J_i`$ on the supergravity side which are conjectured to correspond in the field theory to the voltage and charge under the $`SO(d)`$ R-symmetry group respectively. We can trust the supergravity description of the dual field theory, when the string coupling $`g_s1`$ and the curvatures of the geometry are small. This implies in all cases that the number of coincident $`p`$-branes $`N1`$. For the M2- and M5-brane this is the only requirement. while for the D$`p`$-branes one must further demand that
$$1g_{\mathrm{eff}}^2N^{\frac{4}{7p}},g_{\mathrm{eff}}^2=g_{\mathrm{YM}}^2Nr^{p3}$$
(13)
where for the purpose of the thermodynamics, one needs to set $`r=r_H`$ in the effective YM coupling $`g_{\mathrm{eff}}`$.
In view of this correspondence, we can write the Gibbs free energy and other thermodynamic quantities in terms of field theory variables. In particular we need to specify the relation between the parameter $`\mathrm{}`$ entering the near-horizon limit and the relevant length scale of the theory and compute the rescaled quantities in (2). For the M2 and M5-branes the relation to the 11-dimensional Planck length is $`\mathrm{}=l_p^{3/4}`$ and $`l_p^{3/2}`$ respectively, from which one obtains $`h^6N`$ and $`h^3N`$. For D-branes on the other hand, $`\mathrm{}=l_s`$ and one finds $`Gg_{\mathrm{YM}}^4`$ and $`h^{d2}\lambda `$, where $`g_{\mathrm{YM}}`$ is the Yang-Mills coupling constant and $`\lambda =g_{\mathrm{YM}}^2N`$ the ’t Hooft coupling. Restricting to D$`p`$-branes, the resulting expression for the (low angular momentum expansion of the) free energy (12) in field theory variables is then,
$`F_{\mathrm{D}p}=c_pV_pN^2\lambda ^{\frac{p3}{p5}}T^{\frac{2(7p)}{5p}}[1+{\displaystyle \frac{S_p^1}{\pi ^2}}{\displaystyle \underset{i}{}}\omega _i^2`$
$`+{\displaystyle \frac{S_p^2}{\pi ^4}}\left({\displaystyle \underset{i}{}}\omega _i^2\right)^2+{\displaystyle \frac{S_p^3}{\pi ^4}}{\displaystyle \underset{i}{}}\omega _i^4+\mathrm{}]`$ (14)
where we have defined $`\omega _i=\mathrm{\Omega }_i/T`$ and the numerical coefficients $`c_p`$, $`S_p^a`$ can be found in Ref. . Similar expressions for the M2 and M5 brane, proportional to $`N^{3/2}`$ and $`N^3`$ respectively, are given in this reference as well. The same mapping can be done for all other thermodynamic quantities (9) and we note that for the special value of $`\kappa =\frac{1}{2}(d2)`$ the temperature vanishes, implying that besides the usual extremal limit describing zero temperature field theory, there also exists a limit in which the temperature is zero, accompanied by non-zero R-charges.
## 4 Stability analysis
We can now analyze the critical behavior of the near-horizon limit of the spinning $`p`$-branes, which thus corresponds to the critical behavior of the dual field theories with non-zero voltages under the R-symmetry. The boundaries of stability have been computed in for both the GCE (grand canonical ensemble with thermodynamic variables $`(T,\mathrm{\Omega }_i)`$) and the CE (canonical ensemble, with thermodynamic variables $`(T,J_i)`$) for $`mn`$ equal non-zero angular momenta. To this end, we find the points in phase space at which the functions
GCE: $`det\text{Hes}(F)={\displaystyle \frac{D_{SJ}}{D_{T\mathrm{\Omega }}}}`$ (15)
CE: $`C_J=T\left({\displaystyle \frac{S}{T}}\right)_{\{J_i\}}=T{\displaystyle \frac{D_{SJ}}{D_{TJ}}}`$ (16)
are zero or infinite, for the GCE and CE respectively. Here, we have denoted by $`D_{T\mathrm{\Omega }}`$ the determinant $`\frac{(T,\mathrm{\Omega }_1,\mathrm{\Omega }_2,\mathrm{},\mathrm{\Omega }_n)}{(r_H,l_1,l_2,\mathrm{},l_n)}`$ and likewise for $`D_{TJ}`$, $`D_{S\mathrm{\Omega }}`$ and $`D_{SJ}`$. These determinants can be written as finite polynomials in $`l_i/r_H=J_i/S`$ and their zeroes will determine the boundaries of stability as $`n`$-dimensional submanifolds in the $`(n+1)`$-dimensional phase diagram. Hence they crucially depend on the fact that while there is a one-to-one correspondence between the $`n+1`$ supergravity variables and the extensive quantities $`(S,J_i)`$, the map to the intensive ones $`(T,\mathrm{\Omega }_i)`$ or the mixed combination $`(T,J_i)`$ involves a non-invertible function.
In this talk, we restrict for simplicity to the case of one non-zero angular momentum in the GCE, in which case we find from (15) that the region of stability (for $`d5`$) is determined by the condition
$$J\sqrt{\frac{d2}{d4}}\frac{S}{2\pi }$$
(17)
so that there is an upper bound on the amount of angular momentum the brane can carry in order to be stable. Put another way, at a critical value of the angular momentum density (which equals the R-charge density in the dual field theory) a phase transition occurs. The supergravity description also determines an upper bound on the angular velocity,
$$\mathrm{\Omega }\frac{2\pi }{\sqrt{(d2)(d4)}}T$$
(18)
which is saturated at the critical value of the angular momentum. In Ref. two scenarios have been proposed for the nature of the resulting phase transition: D-brane fragmentation, in which the branes fly apart in the transverse dimension, and phase mixing in which angular momentum localizes on the brane. The latter is possible because from (9) it follows that for each value of $`\mathrm{\Omega }/T`$ between 0 and the maximum in (18), two pairs of supergravity variables $`(r_H,l)`$ can be found, one corresponding to a stable state and the other to an unstable state.
More generally, we remark that for general $`d`$, the two ensembles are not equivalent and that increasing the number of equal-valued angular momenta enlarges the stable region<sup>1</sup><sup>1</sup>1Except for the cases $`d=8,9`$, which have no stability boundary for one non-zero angular momentum.. We also note that the case $`d=4`$ is special since the free energy (12) vanishes, and it is found that the temperature and angular velocity are not independent, so that the phase diagram is degenerate; for example for one non-zero angular momentum one finds the quadratic constraint
$$(2\pi T)^2+\mathrm{\Omega }^2=h^2$$
(19)
Finally, we mention that the critical exponents for all spinning branes can be uniformly computed in both ensembles to be 1/2, a value which satisfies scaling laws in statistical physics.
## 5 Comparison to weak coupling
An important question is to what extent do we observe the above stability phenomena in the large $`N`$ limit of the dual field theory, also at weak coupling. Extending the method of Ref. , one can obtain in an ideal gas approximation the free energies of the field theories for the case of the M-branes and the D-branes of type II string theory,
$$F=\stackrel{~}{c}_pV_pT^{p+1}\underset{\stackrel{}{\alpha }}{}\mathrm{Li}_{p+1}\left[s_\stackrel{}{\alpha }\mathrm{exp}\left(\underset{i=1}{\overset{n}{}}\alpha \omega \right)\right]$$
(20)
where the $`SO(d)`$ R-charge weights $`\stackrel{}{\alpha }`$ run over the 16 different particles and $`s_\stackrel{}{\alpha }`$ is $`+1`$ for bosons and $`1`$ for fermions. The polylogarithms $`\mathrm{Li}_{p+1}`$ are not defined for real numbers greater than one, but can be continued to this region, so that the free energies can be expressed as exact power series in $`\omega _i`$ . From these one can extract the stability behavior by considering again the zeroes of $`det(\text{Hes}(F))`$. The resulting critical values of $`\omega `$ are summarized in Table 5.1, along with the corresponding values (18) at strong coupling.
For zero angular momentum, it has been conjectured that the free energy smoothly interpolates between the weak and strong coupling limit, which is easily extended to the spinning case . In particular, for the D-branes with $`N`$ fixed but with $`\lambda `$ varying between the two limits, the conjecture reads
$$F_\lambda (T,\{\mathrm{\Omega }_i\})=f(\lambda ,T,\{\mathrm{\Omega }_i\})F_{\lambda =0}(T,\{\mathrm{\Omega }_i\})$$
(21)
where $`F_{\lambda =0}(T,\{\mathrm{\Omega }_i\})`$ is the free energy for $`\lambda =0`$. For M-branes the conjectures involves a function interpolating between $`N=1`$ and any finite $`N`$. An important first check of this conjecture is the fact that the free energies of the D-branes in the two limits show the same $`N^2`$ factor in front. This implies that only string loop corrections, which carry factors of $`\frac{1}{N^2}`$ would modify this behavior, and thus do not have to be considered in the large $`N`$ limit.
As a corroboration of this conjecture, it is seen from Table 5.1 that the critical values of the dimensionless quantity $`\mathrm{\Omega }/T`$ for the D2, D3, D4, M2 and M5-branes are remarkably close in the weak and strong coupling limit: It is hence plausible that in these cases the conjectured interpolation between the two limits holds. For the D1 and D6-brane case, however, there cannot be a smooth transition and a boundary of stability is somehow created/destroyed at some special point between the two limits. Finally, for the D5-brane we see that moving away from strong coupling at some point the phase space must expand, since at weak coupling $`T`$ and $`\mathrm{\Omega }_i`$ are again independent thermodynamic quantities.
Turning to the second test, we recall that the string loop expansion in $`g_s`$ and derivative expansion in $`\alpha ^{}=l_s^2`$ translates through the AdS/CFT correspondence into a $`1/N`$ and $`1/\lambda `$ expansion respectively. The tree-level higher derivative term $`l_s^6R^4`$ of the type II string effective action can thus be used to compute the $`\lambda ^{3/2}`$ corrections to the free energy of the D-branes, which was analyzed for the non-rotating D3-brane in . For simplicity we focus on the spinning D3-brane with one non-zero angular momentum, for which the weak and strong coupling free energies are respectively given by,
$`F_{\lambda =0}(T,\mathrm{\Omega })=N^2V_3T^4\left({\displaystyle \frac{\pi ^2}{6}}+{\displaystyle \frac{1}{4}}\omega ^2{\displaystyle \frac{1}{32\pi ^2}}\omega ^4\right)`$
$`F_{\lambda =\mathrm{}}(T,\mathrm{\Omega })`$
$`=N^2V_3T^4\left({\displaystyle \frac{\pi ^2}{8}}+{\displaystyle \frac{1}{8}}\omega ^2+{\displaystyle \frac{1}{16\pi ^2}}\omega ^4+\mathrm{}\right)`$ (22)
Since the dual theory, N=4 $`D=4`$ SYM, is conformal, the interpolating conjecture (21) simplifies to
$$F_\lambda (T,\mathrm{\Omega })=f(\lambda ,\omega )F_{\lambda =0}(T,\mathrm{\Omega })$$
(23)
From (5) $`f(\lambda ,\omega )`$ is thus expected to be smaller than one, decreasing with $`\lambda `$, for fixed $`\omega <1`$, since $`F_{\lambda =\mathrm{}}<F_{\lambda =0}`$ for $`\omega <1`$. This has been corroborated in by computing the $`R^4`$ corrections to the free energy from the spinning D3-brane solution (2).
## 6 Spinning brane bound states and noncommutative field theories
We finally comment on the generalization involving bound states of spinning branes, which are string backgrounds with a non-zero NSNS $`B`$-field (or $`𝒞`$-field in the case of M-theory). These have recently attracted much attention, in view of the appearance of non-commutative geometry in certain limits of such backgrounds, as first discovered in the context of M(atrix) theory . More specifically, non-commutative super Yang-Mills (NCSYM) appears in a special low-energy limit of the world-volume theory of $`N`$ coinciding D$`p`$-branes in the presence of a NSNS $`B`$-field. This fact has been used to extend the correspondence between near-horizon D$`p`$-brane supergravity solutions and super Yang-Mills (SYM) theories in $`p+1`$ dimensions , to a correspondence between near-horizon D$`p`$-brane supergravity solutions with a non-zero NSNS $`B`$-field and NCSYM in $`p+1`$ dimensions .
By applying a sequence of T-dualities on the general spinning D-brane solution (1) one obtains the general spinning D-brane bound state solutions . The resulting backgrounds are bound states of spinning $`D(p2k)`$-branes, $`k=0\mathrm{}m`$, with $`2mp`$ the rank of the NSNS $`B`$-field,
$`B_{2k1,2k}`$ $`=`$ $`\mathrm{tan}\theta _k\left(H^1D_k1\right)`$ (24)
$`D_k`$ $`=`$ $`\left(\mathrm{sin}^2\theta _kH^1+\mathrm{cos}^2\theta _k\right)^1`$ (25)
where $`k=1\mathrm{}m`$ and we refer to for the explicit form of the complete solution. The extra parameters that are introduced are the angles $`\theta _k`$, $`k=1\mathrm{}m`$. Besides the charges and chemical potentials of the D-branes in the bound states, which depend on these angles, the thermodynamic quantities of the bound state solution are not affected by the non-zero $`B`$-field and coincide with those of the spinning D$`p`$-brane given in (9).
The construction of the near-horizon limit of these solutions uses the rescaling (2), and takes the angles $`\theta _k`$ to $`\frac{\pi }{2}`$ while keeping fixed
$$b_k=l_s^2\mathrm{tan}\theta _k$$
(26)
which are the non-commutativity parameters entering the position commutators on the world-volume of the brane
$$[y^{2k1},y^{2k}]b_k,k=1\mathrm{}m$$
(27)
The dual field theories is then NCSYM in $`p+1`$ dimensions with effective coupling constant
$$zg_{\mathrm{eff}}^2=g_{\mathrm{YM}}^2N\left(\underset{k=1}{\overset{m}{}}b_k\right)r^{p3}$$
(28)
It turns out that the gauge coupling phase structure exhibits a far richer structure than the commutative case, as first shown in and generalized in . To study this we set $`b_k=b`$, $`k=1\mathrm{}m`$ and consider a general path in phase space, parametrized by
$$g_{\mathrm{YM}}^2z^\alpha ,rz^\beta ,bz^\gamma $$
(29)
keeping $`N`$ fixed, and where (28) constrains the scaling exponents to satisfy $`\alpha +(p3)\beta +m\gamma =1`$. The case $`\alpha =\gamma =0`$ in (29) corresponds to the one described in . It then follows that the supergravity description is valid provided
$$zN^{\frac{4}{7p}}\left(1+\left(\frac{z}{z_{\mathrm{nc}}}\right)^\eta \right)^{\frac{2m}{7p}}$$
(30)
where we have used that
$$b^2\left(\frac{r}{h}\right)^{7p}=\left(\frac{z}{z_{\mathrm{nc}}}\right)^\eta ,\eta =4\beta +2\gamma 1$$
(31)
Depending on the values of $`\eta `$ and $`z_{\mathrm{nc}}`$ this condition can then be shown to generate four types of phase diagrams, and for each spatial worldvolume dimension of the brane and each non-zero rank of $`B`$-field, a path and region of phase space can be chosen such that the phase structure of any of the four phase diagrams can be realized . These exhibit various interesting features, including regions in which the supergravity description is valid for finite $`N`$ and/or the effective coupling is allowed to range from the transition point $`g_{\mathrm{eff}}1`$ all the way to infinity.
As was argued for the non-rotating case , the thermodynamic quantities are the same as for commutative SYM case up to the replacement
$$g_{\mathrm{YM}}^2g_{\mathrm{YM}}^2\underset{k=1}{\overset{m}{}}b_k$$
(32)
where $`g_{\mathrm{YM}}`$ on the left-hand side is the YM coupling constant of the commutative theory and $`g_{\mathrm{YM}}`$ on the right-hand side of the non-commuta-tive theory<sup>2</sup><sup>2</sup>2For the noncommutative D3-brane case, tree-level $`R^4`$ corrections to the thermodynamics have been addressed in .. This was argued at weak coupling from the field theory point of view by showing that the planar limit of SYM and NCSYM coincide. As an application of the general phase structure analysis, the validity of the thermodynamics for the NCSYM has been examined by requiring that the intensive thermodynamic parameters are invariant for the path in phase space. This condition corresponds to the choice $`\alpha +m\gamma =(5p)\beta `$ in (29). In particular, this enables determining the region of phase space in which $`N`$ can be finite and at the same time the coupling can be taken all the way to infinity. The resulting condition is that
$$Tb^{1/2}\left(\frac{r_0}{r_H}\right)^{\frac{7p}{2}}$$
(33)
showing that at fixed non-extremality parameter $`r_0`$ and horizon radius $`r_H`$, the larger the non-commutativity parameter $`b`$, the larger the temperature region in which these properties are satisfied.
The presence of non-zero angular momenta does not qualitatively change the gauge coupling phase structure, but may well provide further insights into NCSYM in the presence of voltages for the R-charges. Moreover, in view of the recent discovery that the D6-brane theory with $`B`$-field decouples from gravity it is interesting that while the non-rotating case is thermodynamically unstable, for the spinning D6-brane stability is found in the canonical ensemble for sufficiently high angular momentum density . We refer to for a detailed treatment of the thermodynamics of this case.
We finally comment on some issues related to the possibility of finite $`N`$. In the non-commuta-tive case finite $`N`$ is possible when the non-commu-tativity parameter $`a^{\mathrm{eff}}`$ is large, so that the $`1/N`$ expansion<sup>3</sup><sup>3</sup>3For a discussion of $`1/N`$ corrections to the NCSYM thermodynamics, see for an analysis at strong coupling and for a field theory analysis. turns into an expansion in $`1/a^{\mathrm{eff}}`$ in that case . This suggests that for finite $`N`$ and large $`a^{\mathrm{eff}}`$ planar diagrams dominate; indeed this has recently been shown to hold in perturbative field theory . The fact that finite $`N`$ is possible also fits together with the recent observation that there are infinitely many D$`(p2)`$-branes in the near-horizon limit of the D$`(p2)`$-D$`p`$ system, which means that the world volume theory of the D$`(p2)`$-brane has gauge group $`U(\mathrm{})`$.
## Acknowledgments
We thank the organizers of the TMR network meeting Quantum aspects of gauge theories, supersymmetry and unification in Paris for a very pleasant and interesting workshop. |
warning/0002/cond-mat0002382.html | ar5iv | text | # Collective Diffusion and a Random Energy Landscape
## I Introduction
Many systems behave on the phenomenological level essentially randomly and therefore other approaches for the theoretical treatment have to be employed. The randomness, resulting from stochastic forces or be intrinsic in the underlying microscopic theory, inevitably leads to the description of such systems in terms of probabilities and expectation values . The time development of probability is usually found using a master equation. The past years have seen an exciting new development based on the observation of the close relationship between the Markov generator of the master equation and a time evolution operator acting on a many-particle Fock space , for some recent reviews compare . The new insight has led to a series of remarkable exact solutions for the stochastic dynamics of interacting particle systems, for a recent overview see . Despite of exact results the mentioned method has been also fruitful in an approximative description of other models such as the facilitated kinetic Ising system as a candidate for glassy systems or in branching and annihilation random walks . Whereas the original paper , see also , are concerned with a mapping of the master equation to a representation in terms of second-quantized bosonic operators a great progress for exact solvable models had been achieved by mapping to spin-$`1/2`$ Pauli-operators . This mapping to spin systems applies to processes where each lattice site can be occupied by only a finite number of particles. Physically, this restriction may be hard-core constraints or fast on-site annihilation processes. Obviously, such a mapping simulates the exclusion principle for classical lattice models with in a cellular automata.
In the present paper the Fock space description is applied for systems far from equilibrium which are coupled to a heat bath. In particular, we discuss the collective hopping process of a classical many body system, coupled to the mentioned heat bath, and under the influence of a random energy landscape realized by a stochastic activation energy. The particles making random walks have to overcome spatially distributed energy barriers. As the consequence the hopping process is accomplished by a competing force field which can give rise to anomalous diffusion. Further, the analysis should be different considering both cases, the bosonic and the exclusive ones. In the first case the particles should find more rapidly the local energy minima however because of that their mobility could be reduced. As the consequence of the random walk where the particles have to overcome spatially distributed energy barriers, the resulting effective force field can give rise to anomalous diffusion. It is well known that one of the reasons for an anomalous diffusive behavior can be traced back to the influence of a stochastic force field below a critical dimension . An alternative way of self induced anomalous diffusion had been discussed recently introducing a feedback coupling between the diffusive particle and its local environment. Both, the disorder and the memory controlled feedback may lead to a subdiffusive behavior or to localization. The analytical approach could be confirmed by simulations in one and two dimensions where at the critical dimension $`d_c=2`$ logarithmic corrections in the mean square displacement had been found .
Here, we demonstrate that the Fock space approach leads in both cases, bosonic and under exclusion, may lead to anomalous diffusion. Within the long time limit and on a large spatial scale both systems belong to the same universality class.
## II Quantum Approach to Nonequilibrium
The analysis is based on a master equation
$$_tP(\stackrel{}{n},t)=L^{}P(\stackrel{}{n},t)$$
(1)
where $`P(\stackrel{}{n},t)`$ is the probability that a certain configuration characterized by a state vector $`\stackrel{}{n}=(n_1,n_2\mathrm{}n_N)`$ is realized at time $`t`$. There are two special cases, eithe each lattice site is occupied by an arbitrary number of particles $`n_i=0,1,2\mathrm{}`$ or as in a lattice gas $`n_i=0,1`$. Further, the occupation numbers $`n_i`$ are considered as the eigenvalues of the particle number operator defined by creation operators $`d_i^{}`$ or by annihilation operators $`d_i`$. The problem is to formulate the dynamics in such a way that the possible realizations for the occupation numbers are taken into account explicitly. The situation in mind can be analyzed in a seemingly compact form using the master equation in a quantum Hamilton formalism , for a recent reviews see . The dynamics is determined completely by the form of the evolution operator $`L^{}`$, specified below, and the commutation relations of the underlying operators $`d_i^{}`$ and $`d_i`$. Within that approach the probability distribution $`P(\stackrel{}{n},t)`$ is related to a state vector $`F(t)`$ in a Fock-space according to $`P(\stackrel{}{n},t)=\stackrel{}{n}F(t)`$. The basic vectors $`\stackrel{}{n}`$ are composed of the operators $`d_i^{}`$ and $`d_i`$. Using the relation
$$F(t)=\underset{n_i}{}P(\stackrel{}{n},t)\stackrel{}{n}$$
(2)
the master eq. (1) can be transformed into an equivalent one in a Fock-space
$$_tF(t)=LF(t)$$
(3)
where the operator $`L^{}`$ in (1) is mapped onto the operator $`L=\stackrel{}{m}L_{mn}^{}\stackrel{}{n}`$ in eq.(3). It should be emphasized that the procedure is up to now independent on the realization of the basic vectors. Originally, the method had been applied for the Bose case . Recently, an extension to restricted occupation numbers (two discrete orientations) was proposed . Further extensions to p–fold occupation numbers as well as to models with kinetic constraints and to systems with two heat bathes are possible.
As shown by Doi the average of an arbitrary physical quantity $`B(\stackrel{}{n})`$ can be calculated by the average of the corresponding operator $`B(t)`$
$$B(t)=\underset{n_i}{}P(\stackrel{}{n},t)B(\stackrel{}{n})=sBF(t)$$
(4)
with the state function $`s=\stackrel{}{n}`$. The evolution equation for an operator $`B(t)`$ reads now
$$_tB=s[B(t),L]F(t)$$
(5)
As the result of the procedure, all the dynamical equations governed by the classical problem are determined by the structure of the evolution operator $`L`$ and the commutation rules of the operators.
## III Coupling to a Heat Bath
The evolution operator for a collective hopping process is different for an arbitrary occupation number, denoted as Bose case, or an restricted occupation number, denoted as Fermi case. For the last system the operator $`L_f`$ reads
$$L_f=\mu \underset{i,j}{}\left(d_i^{}d_j(1n_i)n_j\right)$$
(6)
where $`\mu `$ is the hopping rate between adjacent sites $`i`$ and $`j`$. The occupation number operator $`n_i=d_i^{}d_i`$ is related to the spin operator by the relation $`S_i=12n_i`$ and the commutation rule is $`[d_i,d_j]=\delta _{ij}(12n_i)`$.
For the Bose case we get
$$L_b=\mu \underset{i,j}{}\left(d_i^{}d_jn_j\right)$$
(7)
where $`d_i^{}`$ and $`d_i`$ fulfills the Bose commutation rules. A generalization to processes under the coupling to a heat bath with a fixed temperature $`T`$ is discussed in . As demonstrated in the evolution operator has to be replaced by
$$L_f=\mu \underset{i,j}{}\left[(1d_id_j^{})\mathrm{exp}(\beta H/2)d_i^{}d_j\mathrm{exp}(\beta H/2)\right]$$
(8)
where the hopping rate $`\mu `$ defines a microscopic time scale; $`\beta =T^1`$ is the inverse temperature of the heat bath and $`H`$ is the Hamiltonian as a measure for the energy. A further generalization is realized by introducing different local heat bathes is discussed in . In the bosonic case the generalization to finite temperatures leads to
$$L_b=\mu \underset{i,j}{}\left[(1\delta _{ij})\mathrm{exp}(\beta H/2)d_i^{}d_j\mathrm{exp}(\beta H/2)\right]$$
(9)
Here we study the case that the Hamiltonian $`H`$ in eqs.(8,9) is simply given by a stochastic energy landscape defined by the energy functional
$$H=\underset{i}{}\epsilon _in_i$$
(10)
Whenever the energy is positive the empty site is energetically favored. Further, $`\epsilon `$ is assumed to be a stochastic local energy the distribution of which will be introduced below based on the continuous representation. In this manner, the model describes a collective hopping process where the jumping particles are subjected to a local random energy $`\epsilon _i`$ which supports or prevents the hopping process with a probability proportional to $`\mathrm{exp}(\pm \epsilon _i/2T)`$. Taking into account the commutation rules we get in both cases
$$e^{\beta H/2}d_ie^{\beta H/2}=d_ie^{\epsilon _i/2T}e^{\beta H/2^{}}d_i^{}e^{\beta H/2}=d_i^{}e^{\epsilon _i/2T}$$
(11)
Using eq.(5) and the algebraic properties of Pauli–operators, the evolution equation for the averaged density reads
$`\mu ^1_tn_r`$ $`=`$ $`{\displaystyle \underset{j(r)}{}}[\mathrm{exp}((\epsilon _j\epsilon _r)/2T)n_j\mathrm{exp}((\epsilon _r\epsilon _j)/2T)n_r`$ (12)
$``$ $`2\mathrm{sinh}({\displaystyle \frac{\epsilon _j\epsilon _r}{2T}})n_rn_j`$ (13)
In the Bose case the evolution equation is much simpler.
$$\mu ^1_tn_r=\underset{j(r)}{}\left[\mathrm{exp}((\epsilon _j\epsilon _r)/2T)n_j\mathrm{exp}((\epsilon _r\epsilon _j)/2T)n_r\right]$$
(14)
Both equations reflect the conservation of the particle number which will be more transparent in a continuum representation. In the special case of a constant energy $`\epsilon _r=\epsilon _j`$ it results the conventional diffusion equation in a discrete version. When the energy changes from site to site the nonlinear eq.(13) is the first step of a whole hierarchy of evolution equations. Assuming now smoothly changing energy $`\epsilon _r`$ and density $`n_r`$ a gradient expansion is appropriate up to the order $`l^2`$ where l is the lattice size. To make the expansion invariant under the underlying rotational symmetry we have to use the following identity
$`{\displaystyle \underset{j(r)}{}}\mathrm{exp}((\epsilon _j\epsilon _r)/2T)n_j`$ $`=`$ $`{\displaystyle \underset{j(r)}{}}n_r`$ (15)
$`+`$ $`\mathrm{exp}(\epsilon _r/2T){\displaystyle \underset{j(r)}{}}\left[\mathrm{exp}(\epsilon _j/2T)n_j\mathrm{exp}(\epsilon _r/2T)n_r\right]`$ (16)
Such an expression reads in a continuous representation including terms of the order $`l^2`$
$$zn(𝐫,t)+\mathrm{exp}(\epsilon (𝐫)/2T)^2[(\mathrm{exp}(\epsilon (𝐫)/2T)n(𝐫,t)]$$
with the averaged density $`n_rn(𝐫,t)`$; $`z`$ is the number of nearest neighbors. After decoupling the nonlinear term in eq.(13) and performing the continuous limit the density $`n(𝐫,t)`$ obeys the following nonlinear diffusion-like equation
$$\mu ^1l^2_tn=^2n+n(1n)\frac{^2\epsilon }{T}+(12n)n\epsilon /T$$
(17)
In a system with exclusion the density couples in a nonlinear manner to the stochastic energy field $`\epsilon (𝐫)`$. Due to the exchange coupling of the evolution operator $`L`$ in eq.(6) the resulting equation (17) is a conserving one where the current is given by
$$𝐣_f=nn(1n)\frac{\epsilon }{T}$$
(18)
In the Bose case we find after performing the continuous limit the density $`n(𝐫,t)`$ obeys the following exact equation
$$\mu ^1l^2_tn=^2n+\frac{1}{T}[n\epsilon ]$$
(19)
The conservation law is manifested in the current
$$𝐣_b=nn\frac{\epsilon }{T}$$
(20)
The resulting equation is nothing else as the conventional diffusion equation under an additional drift term where the Einstein relation is automatically fulfilled. Remark that one can derive a similar equation when the system is coupled to two heat bathes with different temperatures. In that case one has to replace $`\epsilon (𝐫)/T`$ by $`\frac{\nu }{T(𝐫)}`$ where $`\nu `$ is the chemical potential and $`T(𝐫)`$ is the local temperature, see also . In the Bose case eq.(19) depends on the density in a linear manner. It is of Fokker-Planck-type when the density $`n(𝐫,t)`$ is considered as the single probability distribution to find a particle at site $`𝐫`$ at time $`t`$. Such an interpretation is always possible because we have not taken into account any interactions. Therefore, the particles are independent from each other and the concentration field behaves as the probability distribution of a single particle of this system. Different to the case of an arbitrary occupation the current $`𝐣_f`$ includes a term $`n(1n)`$ which is characteristic for systems with exclusion. Due to the exclusion principle the systems reveals a kind of correlation which leads even in the mean field limit to a nonlinear current. Following the discussion for the Bose case eq.(17) can be interpreted as a nonlinear Fokker-Planck-equation for a single particle. The nonlinearity reflects a feedback of a particle to itself due to the excluded volume effect.
It seems to be more appropriate to introduce the force vector $`𝐟(𝐫)=\epsilon (𝐫)`$ the evolution equation in the Bose case reads now
$`\mu ^1l^2n(𝐫,t)`$ $`=`$ $`^2n{\displaystyle \frac{1}{T}}𝐟n{\displaystyle \frac{1}{T}}𝐟n`$ (21)
In the Fermi case the corresponding equation is
$`\mu ^1l^2n(𝐫,t)`$ $`=`$ $`^2n{\displaystyle \frac{1}{T}}𝐟n(12n){\displaystyle \frac{1}{T}}n(1n)𝐟`$ (22)
When the force field $`𝐟(𝐫)`$ is a stochastic one the system offers anomalous diffusive behavior .
## IV Scaling
Now let us discuss both equations when the force field is an stochastic pure spatial dependent field, the correlator of which is given by
$$\overline{f_\alpha (𝐫)f(𝐫^{})}=\varphi _{\alpha \gamma }(𝐫𝐫^{}),\overline{f_\alpha (𝐫)}=0$$
(23)
After averaging over the distribution function of the force field the system is homogeneous depending only on the difference of the spatial coordinates. The most general form of the function $`\varphi _{\alpha \gamma }`$ is given in a Fourier representation by
$$\varphi _{\alpha \gamma }=A(\stackrel{}{q})(\delta _{\alpha \gamma }n_\alpha n_\gamma )+B(\stackrel{}{q})n_\alpha n_\gamma \text{with}n_\alpha =\frac{q_\alpha }{q}$$
(24)
Introducing dimensionless variables $`xx\mathrm{\Lambda }^1,tt\mathrm{\Lambda }^z`$, where z is the dynamical critical exponent and further $`nn\mathrm{\Lambda }^d`$ and according to eq.(24) for constant $`A`$ and $`B`$ $`𝐟𝐟\mathrm{\Lambda }^{d/2}`$ we find the critical dimensionality $`d_c=2`$. For $`d2`$ the term proportional to $`(𝐟n)`$ is relevant whereas the additional term in case of exclusive motion $`n𝐟n`$ is only relevant for $`d<2/3`$. That means for the physical dimension $`d1`$ both models belong to the same universality class, where only $`d2`$ the disorder is relevant. Physically the result is obvious because in the long time limit and for a large spatial scale the Fermi system can be considered to consist of blocks with an increasing size. The larger such a block the more irrelevant is to distinguish both cases, arbitrary occupation and restricted occupation. In case of $`d2`$ the system reveals anomalous diffusive behavior as it had been demonstrated for a similar model not for the density $`n(𝐫,t)`$ but for the probability to $`P`$ to find a particle at time $`t`$ at the point $`𝐫`$. Making the same calculation we end up with the flow equations for the dimensionless coupling parameters $`D=\mu l^2,a=\frac{A}{D^2T^2}K_d,b=\frac{B}{D^2T^2}K_d`$, with $`K_d(2\pi )^d`$: the volume of the d-dimensional unit sphere and $`ϵ=2d`$$`\xi =\mathrm{ln}(\frac{\mathrm{\Lambda }_0}{\mathrm{\Lambda }}`$
$`{\displaystyle \frac{D}{\xi }}`$ $`=`$ $`D\left[z2+{\displaystyle \frac{a(d1)}{d}}{\displaystyle \frac{b}{d}}\right]`$ (25)
$`{\displaystyle \frac{a}{\xi }}`$ $`=`$ $`a\left[ϵa+{\displaystyle \frac{b(d1)}{d}}\right]`$ (26)
$`{\displaystyle \frac{b}{\xi }}`$ $`=`$ $`b\left[ϵ{\displaystyle \frac{a}{d}}\right]`$ (27)
In the same manner one can derive an equation for the mean square displacement $`R=\mathrm{\Lambda }^2s(D,a,b)`$ with $`s=𝐫^\mathrm{𝟐}`$ The flow equation can be written as
$$2s=\frac{s}{D}_\xi D+\frac{s}{a}_\xi a+\frac{s}{b}_\xi b$$
(28)
That equation leads to a scaling behavior of the mean square displacement in the vicinity of the fixed points of eqs.(27). In order to keep the diffusivity $`D`$ fixed to its bare value the effective dynamical exponent $`z(\xi )`$ satisfies $`z(\xi )=2+b(\xi )/d+a(\xi )(1d)/d`$. When the disorder is irrelevant the fixed points are $`a^{}=b^{}=0`$ the exponent is $`z=2`$. For the fixed point $`a^{}=\epsilon d,b^{}=0`$ it results $`z=2\epsilon `$ and for $`a^{}=b^{}=ϵd`$ we find $`z=2+O(ϵ^2)`$. These values are well known . At the critical dimension $`d_c=2`$ we proceed on the following manner. The observation time $`t`$ is related to an initial time $`t_0`$ by
$$t=t_0\mathrm{exp}(_0^\xi z(\xi ^{})𝑑\xi ^{})$$
(29)
Using eqs.(27,28 we can fix the scaling parameter $`\xi `$ according to eq.(29 to be
$$\xi \frac{1}{2}\mathrm{ln}(\frac{t}{t_0})+\frac{1}{2}\mathrm{ln}(1+\frac{a_0}{2}\frac{t}{t_0})$$
where $`a_0`$ is initial value for the parameter $`a`$. From eq.(28) we find the following behavior for the mean square displacement
$$𝐫^2=c_1\frac{t}{t_0}+c_2\frac{t}{t_0}\mathrm{ln}(\frac{t}{t_0})$$
(30)
where $`c_1`$ and $`c_2`$ are two non-universal constants. As expected the system reveals logarithmic corrections at the critical dimension.
## V Behavior above the critical dimension
The thermalized version of the Fock space representation, see eqs.(8,9), leads in the limit $`T\mathrm{}`$ to conventional diffusion. In the high temperature limit the particles are able to overcome each barrier and as the consequence of the stochastic hopping process one finds diffusive behavior in the long time limit independently on the underlying statistics. When the temperature is finite there appears a competition between two processes resulting in a different behavior for both systems. Bose particles can easily find a minimum within the energy landscape defined by the stochastic force. Particles with exclusion have to search for a longer time and on a larger scale to reach an appropriate potential minimum. From here one would conclude to an enhanced diffusivity. On the other hand, the mobility of Bosons is eventually reduced because they find more rapid a stable minimum. Due to the established universality for low dimensions a variation of the behavior should be only observed above the critical dimension. In this regime conventional perturbation theory should be applicable. Let us therefore present lowest order corrections to the the diffusion parameter $`D`$. The effective diffusivity is defined by
$$D_{eff}=\left|\frac{\overline{n^1(\stackrel{}{q},\omega )}}{q^2}\right|_{q=0,\omega =0}$$
(31)
As well as the Bose and the Fermi system lead in second order, proportional to $`\frac{1}{T^2}`$, to nontrivial corrections which are also manifested in the averaged density $`\overline{n(𝐫,t)}`$ or the averaged correlation function $`\overline{n(𝐫,t)n(𝐫^{},t^{})}`$. Indeed, the Fermi system offers additional terms for the density or the correlation function compared with the Bose case. However those terms does not contribute at zero wave vector and hence there are not relevant corrections to the divergent part of $`D_{eff}`$ for $`dd_c`$. Above $`d_c`$ the behavior of the effective diffusion coefficient can be estimated using a perturbative approach around the homogeneous solution denoted by $`\overline{n}`$. We get
$`D_{eff}^f`$ $`=`$ $`D_{eff}^b+{\displaystyle \frac{(1\overline{n})\overline{n}}{DT^2}}I`$ (32)
$`\text{with}I`$ $`=`$ $`{\displaystyle \frac{4K_d}{d}}I_1[BA(d1)]`$ (33)
$`I_1`$ can be expressed by a momentum integral which is always positive in the mesoscopic regime $`\mathrm{\Lambda }>l`$. For $`BA(d1)>0`$, realized for a pure potential field ($`B`$ is the relevant variable, see eq.(24)), eq.(33) leads to
$$D_{eff}^f>D_{eff}^b$$
(34)
Remark that the correction to the bare diffusion coefficient $`D`$ is of the order $`(12\overline{n})^2`$, that means for the half-filled case there are no corrections. That reasonable result should be also valid in a more refined approach.
Because the homogeneous solution is not necessary a stable one we can also estimate the behavior using linear stability analysis around the stationary solution denoted as $`n_s(𝐫)`$. Let us introduce $`n(𝐫,t)=n_s(𝐫)+y(𝐫,t)`$ then the correction $`y(𝐫,t)`$ fulfills in the Bose case the equation
$$_ty=D^2y+\frac{D}{T}(y\epsilon _b)\text{with}𝐟(𝐫)=\epsilon _b(𝐫)$$
(35)
Here $`\epsilon _b(𝐫)=\epsilon (𝐫)v`$ is the true stochastic potential introduced by eq.(10) and $`v`$ plays the role of the chemical potential which regulates the occupation number. In case of the exclusion model the deviation from the stationary solution $`y(𝐫,t)`$ satisfies the same equation however one has to replace the potential in the Bose case, given in eq.(35), by another effective potential
$$\epsilon _b(𝐫)\epsilon _f(𝐫)=2T\mathrm{ln}\left[\frac{\mathrm{cosh}\left(\frac{\epsilon (𝐫)v}{2T}\right)}{\mathrm{cosh}\frac{v}{2T}}\right]$$
(36)
We have gauged the potentials so that for $`\epsilon _f(𝐫)=0`$ also $`\epsilon (𝐫)=0`$. The hopping particles under exclusion are subjected to the modified stochastic energy landscape given by $`\epsilon _f`$. Expanding $`\epsilon _f`$ in terms of $`\epsilon `$ we find the relation
$$\epsilon _f(𝐫)\mathrm{tanh}(\frac{v}{2T})\epsilon (𝐫)$$
(37)
From here it results
$$\overline{\epsilon _f(𝐫)\epsilon _f(0)}\mathrm{tanh}^2(\frac{v}{2T})\overline{\epsilon _b(𝐫)\epsilon _b(0)}$$
(38)
The effective correlator of the disorder in the Fermi system is drastically decreased in comparison to the Bose case. This result is compatible with the previous discussion leading to eqs.(33,34). In particular in the vicinity of half-filling (where the chemical potential $`v`$ is zero) the influence of the disorder is very weak. This special case corresponds to vanishing linear term expanding eq.(36) according to powers of $`\epsilon `$. In the leading order we obtain
$$\epsilon _f(𝐫)\frac{\epsilon ^2(𝐫)}{4T}$$
Different to the Bose case the effective stochastic potential $`\epsilon _f`$, eq.(36), is always positive definite, that means all the deep negative minima of the original stochastic potential become maxima and therefore they are not more available in case of Fermi system. Obviously, they are already occupied and hence they are not accessible for particles.
## VI Conclusions
In the present paper the collective hopping process on a lattice is studied systematically when the particles are subjected to a random energetic landscape manifested by a stochastic energy profile. In particular, we have taken into account both cases, each lattice site is only occupied by one particle or each site can absorb an arbitrary number of particles. Physically, one expects different behavior. Whereas in the situation under exclusion a particle should spend more time for searching an appropriate energy minimum within the stochastic energy the bosons tend to reduce their mobility because they remain for a longer time in the local minima. A further influence on the motion of the particles is given by the coupling to a heat bath which supports the tendency that the system equilibrates. Starting on a master equation in a second quantized form both cases can be easily realized in terms of Bose-operators or spin-$`1/2`$ Pauli-operators. The annihilation and creation process of particles leads in both cases to a density gradient characteristic for a random walk. Due to the additional coupling to a stochastic energy each particle can not follow that gradient simply but it has to overcome an energy barrier at its starting point and at its end point. There appears a conflicting situation that a particle follows the density gradient but the energy at the starting point is higher than at the end point. In this manner it will jump from an occupied to an empty site however under mobilizing a higher amount of energy (lower temperature). The other situation consists of the fact that a particle follows the density gradient and the energy barrier at the starting point is lower than at the end point (high temperature regime). In this case the hopping process is highly supported by the energy landscape whereas in the previous one the process is restricted. As the consequence anomalous diffusive behavior should be realized below the critical dimension.
In the paper we have demonstrated that the Bose-as well as the Fermi-system belong below the critical dimension to the same universality class within the long time limit and on a large spatial scale. For an increasing scale the system can be considered consisting of blocks with an increasing number of particles. Thus, the cases of restricted and unrestricted occupation number per lattice site should be irrelevant. Despite of the universality the density and the correlation function of both systems are different, in particular for an intermediate interval. In particular, we have discussed the situation above the critical dimension where the diffusion constant can offer different behavior in both cases. |
warning/0002/math0002137.html | ar5iv | text | # Cobordism of immersions of surfaces in non-orientable 3-manifolds.
## Introduction
Following the definitions of \[Wel66\], \[Pin85\], \[BS95\] we will say that two immersions $`f`$ and $`f^{}`$ defined on compact closed surfaces not necessarily connected $`F`$ and $`F^{}`$ and taking values in the same 3-manifold $`M`$ are cobordant if there exists a cobordism $`X`$ between $`F`$ and $`F^{}`$ and an immersion $`\mathrm{\Phi }`$ of $`X`$ in $`M\times I`$ that restricts to $`f\times \{0\}`$ and to $`f^{}\times \{1\}`$. Once fixed the manifold $`M`$ the set $`N_2(M)`$ of cobordism classes of immersions of surfaces in $`M`$ is a semi-group with the composition law given by disjoint union. That this be a group when $`M=𝐑^3`$ is given a priori by the fact that inverses are provided by composition with a reflection in a plane. In fact $`N_2(𝐑^3)`$ is isomorphic to $`𝐙/8𝐙`$, as is proved in \[Wel66\]; explicit invariants are given in \[Bro72\] and \[Pin85\]; a generator is the so-called right immersion of Boy of the projective space, see \[Pin85\]. That $`N_2(M)`$ be a group when $`M`$ is a generic manifold is not straightforward; it is proven in \[BS95\] for orientable 3-manifolds that $`N_2(M)`$ is the finite set $`H_2(M,𝐙/2𝐙)\times H_1(M,𝐙/2𝐙)\times 𝐙/8𝐙`$ with a composition law that twists the compositions of the factors. We adapt here their constructions to the non-orientable case, and obtain in theorem 2.1.2 that again $`N_2(M)`$ is a finite group, with support the set $`H_2(M,𝐙/2𝐙)\times H_1(M,𝐙/2𝐙)\times 𝐙/2𝐙`$ and a composition law similar to the one of the orientable case. The crucial point that causes a non-orientable 3-manifold to have a “smaller” group than an orientable 3-manifold with same homology is proposition 2.3.7, and comes from the fact that in a non-orientable environment there exist isotopies that reverse orientation.
The first section is essentially a review of results of \[HH85\], \[BS95\] and \[KT89\]; some remarks and properties are original and are used in the following; in particular in proposition 1.2.1 we compute explicitly the isotropy group for the action of adding kinks introduced in \[HH85\], and in theorem 1.4.4 we classify bands (that is, immersions of annuli or Moebius bands) in a non-orientable 3-manifolds up to regular homotopy. In the second section we develop the computation of the cobordism group.
This paper is an extended version of the talk given in Palermo in September 1999, during the Congress “Proprietà Geometriche delle Varietà Reali e Complesse: Nuovi Contributi Italiani”.
## 1 Some properties of non-orientable 3-manifolds.
From now on a loop $`c`$ in a non-orientable manifold $`X`$ will be said to be orientable in $`X`$ if it preserves the orientation of $`X`$, that is, if $`TX_{|c}`$ is orientable, and will be said to be non-orientable in $`X`$ otherwise.
### 1.1 Projective and anti-equivariant framings.
It is well-known that orientable 3-manifolds are parallelizable. A non-orientable 3-manifold obviously is not, but its tangent bundle is still as simple as a non-orientable vector bundle can be, that is, its structure group can be reduced to $`\{1,1\}`$. A calculus of characteristic classes in \[HH85\] proves in fact that if $`M`$ is a non-orientable 3-manifold then
$$TMdetMdetMdetM.$$
This means that it is possible to define on $`M`$ a projective framing, i.e. a triple of linearly independent vector fields, well-defined up to sign. A triple $`\{v_1,v_2,v_3\}`$ of linearly independent vector fields on the orientation covering space $`\stackrel{~}{M}`$ of $`M`$ is said to be an anti-equivariant framing if
$$v_i(\sigma (P))=\sigma _{}(v_i(P))P\stackrel{~}{M},i=1,\mathrm{},3,$$
where $`\sigma `$ is the covering translation of $`\stackrel{~}{M}`$; an anti-equivariant framing projects on a projective framing. Given a projective framing of $`M`$ one can always find an anti-equivariant framing of $`\stackrel{~}{M}`$ that projects on it: we call such an anti-equivariant framing a lifting of the projective framing.
We will assume, from now on, that all framings are orthonormal with respect to a fixed metric.
### 1.2 Isotropy in adding kinks.
The main result of \[HH85\] is to set up an action $``$ of $`H^1(F,𝐙/2𝐙)`$ on the set $`Imm_\xi (F,M)`$ of regular homotopy classes of immersions in a homotopy class $`\xi `$ of immersions of a surface $`F`$ in a 3-manifold $`M`$. This action is explicitly described in a beautiful geometric way: if $`f`$ is a representative of an element of $`Imm_\xi (F,M)`$ and $`\alpha `$ is in $`H^1(F,𝐙/2𝐙)`$ then a representative of $`f\alpha `$ is obtained by modifying $`f`$ in a tubular neighborhood $`N`$ of a dual curve to $`\alpha `$, in a way the authors describe as adding a kink: in case both of the curve in $`F`$ and its image in $`M`$ are orientable this local modification is the (local) composition of $`f`$ with the immersion of an annulus in a solid torus shown in figure 1; in the other cases the modification is similar.
If $`c`$ is a dual curve to $`\alpha `$ then $`f\alpha `$ will also be denoted by $`fc`$.
The action of adding kinks is proven to be transitive in any homotopy class. Isotropy depends on a property of the class: we say a homotopy class $`\xi `$ is odd if both of $`F`$ and $`M`$ are non orientable and if there exists a self-homotopy $`H`$ of a map $`f\xi `$ such that $`H(x,)`$ is a loop non-orientable in $`M`$ for $`xF`$; we call odd such a homotopy; the action of $`H^1(F,𝐙/2𝐙)`$ on $`Imm_\xi (F,M)`$ has then a group of isotropy of order 2 in any point if $`\xi `$ is odd. If $`\xi `$ is not odd we say it is even, and isotropy is trivial in this case. This means, in particular, that the set $`Imm_\xi (F,M)`$ has the cardinality of $`H^1(F,𝐙/2𝐙)`$ if $`\xi `$ is even, and half its cardinality if $`\xi `$ is odd.
We compute explicitly the isotropy group for the action on odd classes.
###### Proposition 1.2.1
Let $`\xi `$ be a odd homotopy class of immersions of a non-orientable surface $`F`$ in a non-orientable 3-manifold $`M`$. Then the isotropy group of any regular homotopy class in $`\xi `$ is the vector subspace of $`H^1(F,𝐙/2𝐙)`$ generated by $`w_1(F)`$.
Proof. The action of $`H^1(F,𝐙/2𝐙)`$ on $`Imm_\xi (F,M)`$ is not explicitly defined in \[HH85\]. What is stated there is that, given a point $`f`$ of $`Imm_\xi (F,M)`$, it is possible to define a correspondence $`C_f`$ between $`Imm_\xi (F,M)`$ and $`H^1(F,𝐙/2𝐙)`$ which is 1-1 if $`\xi `$ is even and 1-2 if $`\xi `$ is odd; remark that $`C_f^1`$ is then always a function. This action is explicitly defined in \[BS95\] as
$$\begin{array}{cccc}:& Imm_\xi (F,M)\times H^1(F,𝐙/2𝐙)& & Imm_\xi (F,M)\hfill \\ & (f,\alpha )& & \alpha f:=C_f^1(\alpha ).\hfill \end{array}$$
We prove that, if $`\xi `$ is odd, then $`C_f(f)=\{1,w_1(F)\}`$. The correspondence $`C_f`$ is constructed in several steps. The first is to associate to every $`gImm_\xi (F,M)`$ the homotopy class of its differential: that this be a bijection is a consequence of \[Hir59\]; to $`f`$ is then associated $`df`$. The second step is where non-injectivity appears: to each differential one associates two elements of the set $`Bun_f(TF,TM)`$ of homotopy classes of bundle map that commute with $`f`$, via the choice of a projective framing of $`M`$. It follows from the construction given in \[HH85\] that the classes associated to $`df`$ are $`df`$ itself and $`df`$. The following steps give a bijection between $`Bun_f(TF,TM)`$ and $`H^1(F,𝐙/2𝐙)`$, that associates the identity to $`df`$; we call $`w`$ the element of $`H^1(F,𝐙/2𝐙)`$ associated to $`df`$: we are then left to prove that $`w=w_1(F)`$.
Let $`c`$ be a curve on $`F`$. We describe how to determine $`w(c)`$. Consider on $`c`$ the fiber bundle $`\tau :=f^{}(TM)=TF_{|c}\nu _c`$ and the automorphism of $`\tau `$ given by $`1_{TF_{|c}}1_c`$. If $`f(c)`$ is orientable in $`M`$ then the projective framing determines two opposite framings of the trivial bundle $`\tau `$; the map $`1_{TF_{|c}}1_c`$ read in one of these framings gives a closed path in $`SO(3)`$; $`w(c)`$ is then defined to be the class of this path in $`H_1(SO(3),𝐙/2𝐙)=𝐙/2𝐙`$, that does not depend on the choice between the two framings (nor on the choice of $`c`$ between representatives of its class modulo 2) . If $`f(c)`$ is not orientable then $`\tau `$ is not trivial; the projective framing of $`M`$ restricted to $`\tau `$ defines two triples of orthonormal vector fields discontinuous in a point; choose one of them $`\{v_1,v_2,v_3\}`$; define a path of $`SO(3)`$ by
$$P\text{linear transformation between }\{v_1\left(P\right),v_2\left(P\right),v_3\left(P\right)\}\text{ and }1_{TF_{|c}}1_c(v_1\left(P\right),v_2\left(P\right),v_3\left(P\right))\text{,}$$
for all $`Pc`$; thought the vector fields are discontinuous this path is well-defined, continuous and closed, and its homology class is independent of the choice between the two opposite triple of discontinuous vector fields. The value of $`w(c)`$ is this homology class, again considered as an element of $`𝐙/2𝐙`$.
We must now compute these values. Recall that a point of $`SO(3)`$ is identified by a point of $`S^2`$, which indicates the oriented rotation axis, and a number between 0 and $`\pi `$, which indicates the rotation angle taken in the positive sense given by orientation; the facts that the null rotation around any axis is the identity and that two rotations of $`\pi `$ along opposite axes coincide give a bijection between $`SO(3)`$ and the projective space $`𝐏^3`$, which is in fact a well-known homeomorphism. The point in $`SO(3)`$ associated to a point $`P`$ of $`c`$ is the rotation of $`\pi `$ whit axis any of the two normal vectors in $`f(P)`$ to $`f(N)`$, $`N`$ being a tubular neighborhood of $`c`$ in $`F`$. We compute the homology of paths of this sort as the intersection number with a generator of $`H_1(SO(3),𝐙/2𝐙)`$: we take the path of rotations of $`\pi `$ with rotation axis rotating itself of $`\pi `$ in the $`xz`$ plane, from $`(1,0,0)`$ to $`(1,0,0)`$.
Consider first the case that $`f(c)`$ is orientable in $`M`$, that is $`\tau `$ is framed by the projective framing. We then consider a homeomorphism of the standard torus in $`𝐑^3`$ to a tubular neighborhood of $`f(c)`$ which sends the standard framing to the framing of $`\tau `$. The image of $`f(N)`$ results as a band in this standard torus, and the normal vector to this bands is parallel to the $`xz`$ plane (in points where this intersection may be considered transverse) a number of times which is even if the band is orientable and odd if the band is a Moebius band: so we have shown that $`w(c)`$ coincides with $`w_1(F)(c)`$ when $`f(c)`$ is orientable in $`M`$.
We are then left to the case $`f(c)`$ non-orientable in $`M`$. We may suppose to have fixed a lifting to $`\stackrel{~}{M}`$ of the projective framing of $`M`$. Consider the non trivial double covering $`\stackrel{~}{N}`$ of $`N`$, let $`\stackrel{~}{f}`$ be a map of this covering to $`\stackrel{~}{M}`$ that commutes with $`f`$ restricted to $`N`$, let $`\stackrel{~}{\tau }`$ be $`\stackrel{~}{f}^{}(TM)`$: this is a trivial bundle framed by the fixed framing of $`\stackrel{~}{M}`$. Again we fix a homeomorphism of the standard torus with standard framing to a tubular neighborhood of $`\stackrel{~}{f}(\stackrel{~}{c})`$ in $`\stackrel{~}{M}`$ with the fixed framing; the image of the band $`\stackrel{~}{f}(\stackrel{~}{N})`$ in this torus gives count two times of the band $`f(N)`$, so consider its (transverse) intersection with, say, the semi-space $`y0`$; the number we are looking for is the number of times, modulo 2, that the normal vector to this part of the band is parallel to the $`xz`$ plane (again this intersection can be taken to be transverse); again this number is even if $`N`$ is orientable and odd if $`N`$ is a Moebius band, so that again $`w(c)=w_1(F)(c)`$, and this ends the proof. $`\mathrm{}`$
### 1.3 Bands in orientable 3-manifolds.
The meaning of the main theorem in \[HH85\] is that the regular homotopy class of a map in its homotopy class is determined by the behaviour of the map in tubular neighborhoods of the curves of the surface. These tubular neighborhoods are either annuli or Moebius bands; we call band in a 3-manifold the immersion of an annulus or a Moebius band. In \[BS95\] the rôle of bands in this subject becomes more evident, and their properties are also a necessary brick in the computation of cobordism in the case of orientable 3-manifolds.
We review the properties of bands in orientable 3-manifolds, then we look at the case of bands in a non-orientable 3-manifold.
#### 1.3.1 Bands in $`𝐑^3`$.
Consider an embedded band in $`𝐑^3`$, and consider the linking number of its core with its (possibly non-connected) boundary. We think of this linking number in the following way: we pick on the boundary of a tubular neighborhood $`N`$ of the core of the band a preferred basis for $`H_1(N,𝐙)`$, with the orientation of the longitude coherent with an orientation of the core; we then give to the meridian the orientation that makes the global orientation off $`N`$ compatible with the orientation of the environment; the intersection of the band with $`N`$ represents a homology class, and the linking number is the meridian coordinate of this class in the preferred basis. This number is even if the band is orientable and odd if it is a Moebius band, so that its class modulo 2 is a total topological invariant of the band.
Given any band in $`𝐑^3`$ there exists a regular homotopy between its core and the standard circle in the $`xy`$ plane in $`𝐑^3`$; extend it to a regular homotopy between a tubular neighborhood of the core of the band and the standard torus of $`𝐑^3`$: we obtain a band regularly homotopic to the original and that differs from the annulus in the $`xy`$ plane by a number of half twists given by the linking number. Now if we isolate 4 half twists there exists a regular homotopy relative to the rest of the band which transforms 4 half twists in a couple of kinks and then in a piece of band with no twisting, see figures 2 and 3;
this local construction reduces by regular homotopy the original band to one of four models, classified by the linking number modulo 4. This number is in fact an invariant in $`𝐙/4𝐙`$ of the regular homotopy class of the band (see \[GM86\], page 114), so it coincides with the original linking number (modulo 4). This proves the following:
###### Proposition 1.3.1
In $`𝐑^3`$ the linking number modulo 4 between the core of a band and the boundary oriented in a coherent way is a total invariant of the regular homotopy class of the band. $`\mathrm{}`$
We call this invariant in $`𝐙/4𝐙`$ number of half twists of the band.
#### 1.3.2 $`Spin`$ structures and preferred longitudes in orientable 3-manifolds.
The problem in extending the definition of half twists of a band to generic 3-manifolds is that in a 3-manifold there is no preferred basis for the homology of the boundary of the tubular neighborhood of a knot: a meridian is still defined, and the orientation of the environment can give it an orientation as it does in $`𝐑^3`$, but there is no way, a priori, to distinguish between longitudes. Remark that the local modification of figure 3 can be performed in any manifold, so that up to regular homotopy we are still interested in the class of rest modulo 4 of the integer coordinate; this means that it is enough to choose a homology class of longitudes modulo 2, since longitudes that differ by two meridians give integer coordinates belonging to the same class modulo 4 (remark that the boundary of a band covers the core twice). In \[KT89\] it is shown how to use a $`Spin`$ structure on an orientable 3-manifold to determine this choice, that reduces to the classical one when $`M=𝐑^3`$ with its unique $`Spin`$ structure. We give here a slightly different way to define the same choice.
Orientable vector bundles on a circle are always trivial. Framings of an oriented vector bundle of rank 3 on a circle are 2, up to homotopy. This is because a homotopy class of framings can be considered as a homotopy class of sections of the associated principal bundle, which, being trivial, is homeomorphic to $`S^1\times SO(3)`$: and $`\pi _1(S^1\times SO(3))`$ contains four classes, only two of which can be realized as sections of the bundle.
Consider $`S^1`$ embedded in $`𝐑^3`$ in the standard way, and consider on it the fiber bundle $`\tau `$ given by restriction of the tangent bundle to $`𝐑^3`$. Define on $`\tau `$ two framings containing in each point the tangent vector to $`S^1`$: let $`e_0`$ be the oriented framing containing the vector parallel to the $`z`$-axis as second element; let $`e_1`$ be the framing containing in the point $`(\mathrm{cos}\theta ,\mathrm{sin}\theta ,0)`$ of $`S^1`$ the vector $`\mathrm{cos}\theta (0,0,1)+\mathrm{sin}\theta (\mathrm{cos}\theta ,\mathrm{sin}\theta ,0)`$ as second element. These two framings are not equivalent, since the curves they describe in the principal bundle of $`\tau `$ differ by a generator of the homotopy group of $`SO(3)`$: this means that any other framing of $`\tau `$ is equivalent either to $`e_0`$ or to $`e_1`$. In particular the standard framing is homotopic to $`e_1`$, since $`e_1`$ extends to a framing of the tangent bundle to $`𝐑^3`$ restricted to the disc, by
$$(\rho \mathrm{cos}\theta ,\rho \mathrm{sin}\theta ,0)((\mathrm{sin}\rho \theta ,\mathrm{cos}\rho \theta ,0),(\mathrm{sin}\rho \theta \mathrm{cos}\rho \theta ,\mathrm{sin}^2\rho \theta ,\mathrm{cos}\rho \theta ),\text{third orthonormal}),$$
and remark that the framing of the trivial bundle on the disc is unique up to homotopy, the disc being contractile. But now the second vector of either $`e_0`$ or $`e_1`$ frames the normal bundle to $`S^1`$ in $`𝐑^3`$, and in particular detects a longitude; so we get close to what we are looking for. Remark that the longitude which is preferred in the classical definition (that is, the one having linking number 0 with the core) is the one given by $`e_0`$, and that the longitude detected by $`e_1`$ belongs to the other class.
On the other side, $`Spin`$ structures on an oriented vector bundle on a manifold $`B`$ are acted on simply transitively by $`H^1(B,𝐙/2𝐙)`$, so that $`Spin`$ structures on an oriented vector bundle of rank 3 on the circle are again 2. It is also true that a framing of a trivial bundle induces naturally a $`Spin`$ structure: in fact a $`Spin`$ structure on a vector bundle can be defined as a double covering of the associated principal $`SO(n)`$-bundle which be non trivial when restricted to the fiber (see \[Mil62\]), and a framing, giving a bundle equivalence between the principal bundle and $`SO(n)\times B`$, induces such a covering by pull-back of the standard covering by $`Spin(n)\times B`$:
The definition of $`Spin`$ structure via double covering implies that $`Spin`$ structures can be described by the elements of $`H^1(\text{principal bundle},𝐙/2𝐙)`$ that restrict to the fiber as generators of $`H^1(SO(n),𝐙/2𝐙)`$ . Two $`Spin`$ structures are equivalent if the two double coverings are equivalent, i. e. if the associated cohomology classes are the same.
Now, $`\pi _1(S^1\times SO(3))`$ and $`H^1(S^1\times SO(3),𝐙/2𝐙)`$ are isomorphic, and the correspondence we gave between homotopy classes of sections of the principal bundle and cohomology classes that restrict to a fiber as a generator is bijective in this case, so that speaking of $`Spin`$ structures or of framings of an oriented vector bundle of rank 3 on $`S^1`$ is the same thing.
We are now ready for the definition. Let $`M`$ be an oriented 3-manifold with a fixed $`Spin`$ structure, let $`K`$ be a knot in $`M`$; the $`Spin`$ structure restricted on $`TM_{|K}`$ gives a framing; choose a diffeomorphism of the standard $`S^1`$ in $`𝐑^3`$ with standard framing to $`K`$ in $`M`$ with this framing, deform by a homotopy the standard framing to $`e_1`$ and pull this deformation back to $`M`$; the second vector of the framing, with usual identification of the normal bundle with the tubular neighborhood $`N`$ of the knot, gives now a curve on $`N`$; the preferred longitude on $`N`$ will be the homology class in $`H_1(N,𝐙/2𝐙)`$ that doesn’t contain this curve. It is straightforward to verify that this choice doesn’t depend on the diffeomorphism nor on the deformation: in fact an oriented framing of $`T𝐑^3`$ restricted to the standard $`S^1`$ and containing the tangent vector to $`S^1`$ is homotopic either to $`e_0`$ or to $`e_1`$ according to the homology class modulo 2 of the curve they pick on the standard torus is equal to the class to the standard longitude or not. This proves the following characterization:
###### Proposition 1.3.2
Any deformation of the framing of $`TM_{|K}`$ that takes it to a framing containing the tangent vector to $`K`$ chooses a curve in $`N`$ which is the non-preferred longitude. $`\mathrm{}`$
This description of the preferred longitude is particularly useful when the $`Spin`$ structure of the manifold is itself associated to a framing of the whole manifold. In this case to assign the preferred longitude it is even not necessary to mention the $`Spin`$ structure.
We prove that our definition coincide with the definition of even framing in \[KT89\], page 209:
###### Lemma 1.3.3
Let $`M`$ be an oriented 3-manifold with a fixed $`Spin`$ structure, and let $`K`$ be a knot in $`M`$; then any even framing of the normal bundle to $`K`$ represents a preferred longitude.
Proof. The choice of even framing only depends on the $`Spin`$ structure of $`TM`$ restricted to $`K`$, so it coincides either with the choice of preferred longitude or with its opposite. But the two choices coincide in the case of the standard circle in $`𝐑^3`$, where they both assign the class of the standard longitude of the torus, so they coincide in every case. $`\mathrm{}`$
If the knot represents the trivial class in $`H_1(M,𝐙/2𝐙)`$ the preferred longitude does not depend on the choice of the $`Spin`$ structure, since any two structures differ by the action of an element of $`H^1(M,𝐙/2𝐙)`$, that doesn’t affect this knot. In this case the preferred longitude can be realized in a geometric way:
###### Proposition 1.3.4
If a knot $`K`$ represents the trivial class in $`H_1(M,𝐙/2𝐙)`$ then take any embedded surface $`F`$ such that $`F=K`$; the (transverse) intersection of $`F`$ and $`N`$ is a preferred longitude.
Proof. This property is the content of theorem 4.3 in \[KT89\]$`\mathrm{}`$
###### Corollary 1.3.5
If $`M=𝐑^3`$ the choice of preferred longitude reduces to the classical. $`\mathrm{}`$
Remark that the definition of preferred longitude is invariant under regular homotopy. For more details on this subject see \[KT89\].
In \[BS95\] the definition of even longitude allows to extend the notion of number of half twists to bands in generic oriented 3-manifolds. What is (implicitly) obtained is that given a knot $`K`$ there are 4 regular homotopy classes of bands having $`K`$ (or a knot homotopic to it) as core, and that the number of half twists is a total invariant.
### 1.4 Bands in non-orientable 3-manifolds.
#### 1.4.1 Equivariant $`Spin`$ structure.
The first and less expensive attempt to extend the definition of half twists to bands immersed in a non-orientable 3-manifold $`M`$ is to look only at bands whose core is orientable in $`M`$, that is, bands that have two homeomorphic preimages in the orientation double covering $`\stackrel{~}{M}`$ of $`M`$: the number of half twist of the band can be defined as the number of half twists of one of its preimages, provided we fix on $`\stackrel{~}{M}`$ a $`Spin`$ structure that makes the choice between the two preimages inessential. Call $`\sigma `$ the covering translation of $`\stackrel{~}{M}`$: we say a $`Spin`$ structure on $`\stackrel{~}{M}`$ is equivariant if, given an embedded circle $`K`$ in $`\stackrel{~}{M}`$ such that $`\sigma KK=\mathrm{}`$, whenever $`l`$ is a curve on the boundary of a tubular neighborhood of $`K`$ that represents a preferred longitude, then $`\sigma l`$ represents a preferred longitude.
###### Lemma 1.4.1
The $`Spin`$ structure on $`\stackrel{~}{M}`$ induced by an anti-equivariant framing is equivariant.
Proof. Remark that $`\sigma `$ restricted to a tubular neighborhood $`N`$ of $`K`$ does not preserve the homotopy class of the framing, since $`\sigma `$ is orientation reversing. Call $`\{v_1,v_2,v_3\}`$ the given anti-equivariant framing: the longitude $`l`$ belongs to the class which doesn’t contain the curve detected on $`N`$ by a deformation of $`\{v_1,v_2,v_3\}_{|K}`$; the same deformation composed with $`\sigma `$ gives a deformation of $`\{\sigma v_1,\sigma v_2,\sigma v_3\}_{|\sigma K}=\{v_1,v_2,v_3\}_{|\sigma K}`$, which picks on the boundary of the tubular neighborhood of $`\sigma K`$ the longitude $`\sigma l`$.
We are then reduced to show that the curve on the standard torus in $`𝐑^3`$ detected by the non-oriented framing $`e_1`$ belongs to the class which doesn’t contain the standard longitude (recall we are talking of classes modulo 2). But this is evident. $`\mathrm{}`$
Now put on $`\stackrel{~}{M}`$ an equivariant $`Spin`$ structure. Let $`\mathrm{\Sigma }`$ be a band in $`M`$ whose core is orientable in $`M`$, let $`\stackrel{~}{\mathrm{\Sigma }}`$ be one of its liftings to $`\stackrel{~}{M}`$; the other is $`\sigma `$$`\stackrel{~}{\mathrm{\Sigma }}`$. We have
###### Proposition 1.4.2
The number of half twists of $`\stackrel{~}{\mathrm{\Sigma }}`$ is opposite to the number of half twists of $`\sigma `$$`\stackrel{~}{\mathrm{\Sigma }}`$.
Proof. Call $`K`$ the core of $`\mathrm{\Sigma }`$, with a fixed orientation, and call $`N`$ a tubular neighborhood of $`K`$, see figure 4.
Let $`\{m,l\}`$ be a preferred basis for $`H_1(N,𝐙)`$: recall the orientation of $`l`$ is chosen according to the orientation of $`K`$, and the orientation of $`m`$ is the one that, together with the orientation of $`l`$, gives to $`N`$ the orientation coming from $`\stackrel{~}{M}`$. This choice makes the meridian coordinate of a curve in $`N`$ independent of the choice of orientation for $`K`$.
Recall that $`\sigma `$ is orientation reversing. This means that $`\{\sigma m,\sigma l\}`$ is not a preferred basis for the homology of a tubular neighborhood of $`\sigma K`$, whereas $`\{\sigma m,\sigma l\}`$ is one. This gives the thesis. $`\mathrm{}`$
But an even number and its opposite are congruent modulo 4; this proves:
###### Corollary 1.4.3
The number of half twists is well defined for orientable bands whose core is orientable in $`M`$$`\mathrm{}`$
#### 1.4.2 Regular homotopies between bands in non-orientable 3-manifolds.
In trying to extend further the definition of number of half twists to generic bands in non-orientable 3-manifold we look at regular homotopies in a non-orientable environment: we will see that in a non-orientable environment there are “more” regular homotopies. First recall the notion of odd self-homotopy in a non-orientable manifold $`X`$ (see § 1.2): a self-homotopy $`H`$ is odd if the closed path $`H(x,)`$ is non-orientable in $`X`$ for any $`x`$.
###### Theorem 1.4.4
Let $`M`$ be a non-orientable 3-manifold, and let $`K`$ be an embedded circle in $`M`$; then
1. if $`K`$ is orientable in $`M`$ and doesn’t admit any odd self-homotopy then there are 4 regular homotopy classes of bands having a circle homotopic to $`K`$ as core;
2. if $`K`$ is orientable in $`M`$ and admits an odd self-homotopy then there are 3 regular homotopy classes of bands having a circle homotopic to $`K`$ as core, in particular any two non-orientable bands of this type are regularly homotopic;
3. if $`K`$ is non-orientable in $`M`$ then there are 3 regular homotopy classes of bands having a circle homotopic to $`K`$ as core, in particular any two non-orientable bands of this type are regularly homotopic; the two classes of orientable bands differ by a reparametrization.
Proof. Let $`\mathrm{\Sigma }`$ be a band whose core is homotopic to $`K`$. Take the core of $`\mathrm{\Sigma }`$ to $`K`$ with a regular homotopy and extend it (for example via an exponential map) to a tubular neighborhood: this takes $`\mathrm{\Sigma }`$ to a band regularly homotopic and having $`K`$ as core. We now think of $`\mathrm{\Sigma }`$ as a band having $`K`$ as core. By means of the local modification shown in figure 3 we can reduce any band having $`K`$ as core to one of four models; call $`\mathrm{\Sigma }_0`$ one of the models of orientable band, the others differ from it only locally: call adding a local twist the operation of substituting to a piece looking as in the left side of figure 5
a piece looking as in the right side of the same figure; $`\mathrm{\Sigma }_i`$ has $`i`$ local half twists, $`i`$ being 1 or 2, $`\mathrm{\Sigma }_1`$ has one local half twist in the opposite direction to the half twist of $`\mathrm{\Sigma }_1`$. We have to decide under which conditions these models can or cannot admit regular homotopies. Recall that bands with number of half twists that differ modulo 2 are defined on non-homeomorphic domains, so that the notion of regular homotopy is only possible between $`\mathrm{\Sigma }_0`$ and $`\mathrm{\Sigma }_2`$ and between $`\mathrm{\Sigma }_1`$ and $`\mathrm{\Sigma }_1`$.
Let now $`K`$ be a curve orientable in $`M`$. Its tubular neighborhood $`N`$ is then a solid torus. The embedding $`K`$ lifts to $`\stackrel{~}{M}`$ in two embeddings $`\stackrel{~}{K}`$ and $`\sigma `$$`\stackrel{~}{K}`$; any homotopy between the two liftings projects to an odd self-homotopy of $`K`$, any self-homotopy of one of the liftings projects to an even self-homotopy of $`K`$. The liftings of the embedding of $`N`$ are two solid tori, and the covering projection restricted to either of these tori gives opposite orientations to $`N`$, see figure 4. If we extend a homotopy between $`\stackrel{~}{K}`$ and $`\sigma `$$`\stackrel{~}{K}`$ to tubular neighborhoods it projects to an odd self-homotopy of $`N`$ that reverses the orientation, if we extend a self-homotopy of either $`\stackrel{~}{K}`$ or $`\sigma `$$`\stackrel{~}{K}`$ to a tubular neighborhood it projects to an even self-homotopy that preserves the orientation. On the other side any self-homotopy of $`K`$ lifts either to a self-homotopy of $`\stackrel{~}{K}`$, if it is even, or to a homotopy between $`\stackrel{~}{K}`$ and $`\sigma `$$`\stackrel{~}{K}`$, if it is odd.
Now if $`K`$ does not admit any self-homotopy (for example if $`\stackrel{~}{K}`$ and $`\sigma `$$`\stackrel{~}{K}`$ are not homotopic) then any regular homotopy between two different models would lift to a regular homotopy between bands having $`\stackrel{~}{K}`$ as core and different number of half twists, and this is not possible by the classification of bands in orientable manifolds: this proves the first part of the theorem.
If $`K`$ admits an odd self-homotopy then take a regular odd self-homotopy and extend it to a regular self-homotopy $`H`$ of $`N`$. Call $`\stackrel{~}{\mathrm{\Sigma }}_i`$ the lifting of $`\mathrm{\Sigma }_i`$ to a band having $`\stackrel{~}{K}`$ as core; we might suppose that $`\stackrel{~}{\mathrm{\Sigma }}_i`$ has $`i`$ half twists. The lifted regular homotopy $`\stackrel{~}{H}`$ takes $`\stackrel{~}{\mathrm{\Sigma }}_i`$ to a band having $`\sigma `$$`\stackrel{~}{K}`$ as core, and the same number of half twists. Recall that, by 1.4.2, $`\sigma \stackrel{~}{\mathrm{\Sigma }}_i`$ has $`i`$ half twists, so, if $`i`$ is odd it is not the final image of $`\stackrel{~}{H}`$. But $`\sigma \stackrel{~}{\mathrm{\Sigma }}_i`$ projects to $`\mathrm{\Sigma }_i`$, so that the final image of $`\mathrm{\Sigma }_1`$ under $`H`$ is the other model, that is, $`\mathrm{\Sigma }_1`$. The same construction applied to $`\mathrm{\Sigma }_0`$ shows that a odd homotopy preserves the model when the band is orientable, and that even homotopy preserve comes from the first part. This proves the second assertion.
We are left to the case when $`K`$ is a non-orientable curve in $`M`$. Its tubular neighborhood $`N`$ is then a solid Klein bottle. We can fix a diffeomorphism of $`N`$ to the following model of solid Klein bottle: consider $`D^2\times I`$ modulo the relation $`(P,0)(\rho (P),1)`$, $`\rho `$ being the reflection of $`𝐑^2`$ in the $`x`$ axis restricted to $`D^2`$. We can consider that $`\mathrm{\Sigma }_0`$ is $`[1,1]\times \{0\}\times I/`$. Remark now that, when we rotate $`[1,1]\times \{0\}\times \{0\}`$ in the positive direction, the other copy of this segment, that is $`[1,1]\times \{0\}\times \{1\}`$, must rotate in the negative direction, see figure 6.
The effect on the band is that, after a rotation of $`\pi `$, $`\mathrm{\Sigma }_0`$ becomes a band having the same support as $`\mathrm{\Sigma }_2`$, and differing from it only by a reparametrization. Remark that the two boundary components of $`\mathrm{\Sigma }_0`$ belong to different homology classes.
On the other side if we apply this isotopy to $`\mathrm{\Sigma }_1`$ and rotate of $`\pi /2`$ we obtain a band with support $`\{0\}\times [1,1]\times I/`$; if we apply the opposite rotation to $`\mathrm{\Sigma }_1`$ we obtain a band with the same support: but this time sliding the band across the orientation disk we obtain exactly $`\mathrm{\Sigma }_1`$, and this ends the proof. $`\mathrm{}`$
This theorem shows that it is not possible to extend the definition of number of half twists to generic bands in non-orientable $`M`$, since they can be equivalent up to regular homotopy plus reparametrization or simply up to regular homotopy.
## 2 Computation of cobordism.
We adapt this proof to the non-orientable case.
### 2.1 Statement and scheme of proof.
Consider an immersion $`f`$ of a compact closed surface $`F`$, non necessarily connected nor orientable, in a non-orientable connected 3-manifold non necessarily compact nor closed. Up to regular homotopy $`f`$ can be considered generic: in this case this means it has a finite number of curves of double points and a finite number of triple points. We define 3 invariants associated to $`f`$: let $`H_f`$ be the homology class represented by $`f`$ in $`H_2(M,𝐙/2𝐙)`$; let $`\delta _f`$ be the homology class represented by the locus of double points of $`f`$, that is, the points of $`M`$ having 2 or 3 preimages, in $`H_1(M,𝐙/2𝐙)`$; let $`n_f`$ be the Euler characteristic of $`F`$ modulo 2 (if $`F`$ is not connected we consider the sum of the Euler characteristics of its connected components).
###### Lemma 2.1.1
These functions are invariant up to cobordism.
Proof. Homology between cobordant immersions is provided by the immersion of the cobordism, so that $`H_{}`$ is well-defined up to cobordism. Homology between the loci of double points of cobordant immersions is provided by the locus of double points of the immersion of the cobordism, so that $`\delta _{}`$ is well-defined up to cobordism. As long as $`n_{}`$ is concerned, recall that the Euler class modulo 2 is the total invariant of classical cobordism of surfaces, so that in particular it is invariant up to cobordism of immersions. $`\mathrm{}`$
If $`Y`$ is a cobordism class of immersions we can then define $`H_Y`$, $`\delta _Y`$ and $`n_Y`$ as the invariants associated to any generic representative of $`Y`$. We will prove that the function
$$\begin{array}{cccc}\hfill \mathrm{\Psi }:& N_2(M)& & H_2(M,𝐙/2𝐙)\times H_1(M,𝐙/2𝐙)\times 𝐙/2𝐙\\ & Y& & (H_Y,\delta _Y,n_Y)\end{array}$$
is a total invariant. First remark that
$$\mathrm{\Psi }(Y+Y^{})=(H_Y+H_Y^{},\delta _Y+\delta _Y^{}+H_YH_Y^{},n_Y+n_Y^{}),$$
where $``$ is the bilinear form of intersection; it is straightforward to check that the composition law
$$(H,\delta ,n)(H^{},\delta ^{},n^{})=(H+H^{},\delta +\delta ^{}+HH^{},n+n^{})$$
makes $`H_2(M,𝐙/2𝐙)\times H_1(M,𝐙/2𝐙)\times 𝐙/2𝐙`$ a commutative group, and $`\mathrm{\Psi }`$ is clearly an homomorphism for this group structure.
###### Theorem 2.1.2
$`\mathrm{\Psi }`$ is an isomorphism of groups between $`N_2(M)`$ and $`(H_2(M,𝐙/2𝐙)\times H_1(M,𝐙/2𝐙)\times 𝐙/2𝐙,)`$.
Proof. We are left to prove that $`\mathrm{\Psi }`$ is bijective.
We first see surjectivity. Given a triple $`(H,\delta ,n)`$ consider: an embedding $`f`$ that represents $`H`$; an embedded circle $`K`$ representing $`\delta `$, its tubular neighborhood $`N`$ and the immersion $`h`$ obtained by adding a kink along any longitude of the inclusion of $`N`$ (remark this is the immersion of a torus or of a Klein bottle according to $`w_1(M)(\delta _Y)`$ being 0 or 1); if $`n_f=n`$ then $`\mathrm{\Psi }(f+h)=(H,\delta ,n)`$; if $`n_fn`$ then consider $`g=\varphi \gamma `$, $`\gamma `$ being a generator of $`N_2(𝐑^3)`$ and $`\varphi `$ a diffeomorphism of $`𝐑^3`$ to a ball $`B`$ in $`M`$: then $`\mathrm{\Phi }(f+h+g)=(H,\delta ,n)`$; and this proves surjectivity.
Consider an immersion $`f`$ of $`F`$ in $`M`$ and let $`l`$ be a closed circle in $`F`$; we denote by $`q_f(l)`$ the number of half twists of the band given by the restriction of $`f`$ to a tubular neighborhood of $`l`$ in $`F`$, possibly perturbed by a regular homotopy. Remark that, since $`M`$ is non-orientable, this is only defined when $`l`$ is orientable in $`F`$ and $`f(c)`$ is orientable in $`M`$. To prove injectivity we show in lemma 2.2.1, adapting the argument in \[BS95\], that any generic representative of a class $`Y`$ can be deformed via surgeries to a representative $`f+h+g`$ such that: $`f`$ is an embedding; $`h`$ is the immersion obtained from the inclusion of the boundary of a tubular neighborhood of a curve representing $`\delta _Y`$ by adding a kink along a longitude, possibly such that $`q_h(l)=0`$ if $`\delta _Y`$ is orientable; $`g`$ is an immersion contained in a ball in $`M`$. This decomposition splits $`\mathrm{\Psi }`$ in its 3 components, that is $`\mathrm{\Psi }(Y)=(H_f,\delta _h,n_f+n_g)`$, so that if $`\mathrm{\Psi }(Y)=\mathrm{\Psi }(Y^{})`$ we are left to show that
$`H_f=H_f^{}`$ $`f_cf^{}`$ (1)
$`\delta _h=\delta _h^{}`$ $`h_ch^{}`$ (2)
$`n_g=n_g^{}`$ $`g_cg^{},`$ (3)
where $`_c`$ means cobordism relation. That (1) is true is the content of lemma 2.3.1; that (2) is true is the content of lemmas 2.3.4 and 2.3.6; that (3) is true is the content of lemma 2.3.8$`\mathrm{}`$
### 2.2 Decomposition lemma.
###### Lemma 2.2.1
Any generic representative of a class $`YN_2(M)`$ can be deformed via surgeries to a representative $`f+h+g`$ such that: $`f`$ is an embedding; $`h`$ is the immersion obtained from the inclusion of the boundary of a tubular neighborhood of a curve representing $`\delta _Y`$ by adding a kink along a longitude, possibly such that $`q_h(l)=0`$ if $`\delta _Y`$ is orientable; $`g`$ is an immersion contained in a ball in $`M`$.
Proof. Let $`Y`$ be a cobordism class, and let $`f_1`$ be a generic representative of $`Y`$: it is an immersion with a finite number of curves of double points and a finite number of triple points. We call $`C(M)`$ the subgroup of $`N_2(M)`$ given by classes which admit a representative immersed in a disk of $`M`$.
First we eliminate triple points. To do so call $`\gamma `$ the right immersion of Boy in $`𝐑^3`$; recall $`\gamma `$ is an immersion with a single triple point; consider the connected sum of $`f_1`$ and $`\gamma `$ in a chart of $`M`$; it is possible to deform by regular homotopy this immersion to an immersion with one triple point less (remark that $`\gamma `$ in a chart of $`M`$ is cobordant to its inverse, see proposition 2.3.7). By recursively repeating this construction we obtain an immersion $`f_2`$ without triple points and cobordant to $`f_1`$ up to an element of $`C(M)`$.
The immersion $`f_2`$ has a finite number of curves of double points. Take two such curves, connect them with a path, then substitute the path with a couple of tubes as in figure 7.
By recursively repeating this operation we are then left with an immersion $`f_3`$ representing $`Y`$ up to an element of $`C(M)`$ and with a single curve $`K`$ of double points; clearly $`K`$ represents $`\delta _Y`$.
The intersection of the image of $`f_3`$ with a tubular neighborhood of $`K`$ is a bundle on $`S^1`$ with fiber a figure $`X`$, orientable or non-orientable according to the orientability of $`\delta _Y`$. If we number orderly the edges of the figure $`X`$ we can identify the its group of isometries with a subgroup of $`𝒮_4`$; in this framework we can reduce the structure group of our fiber bundle to one and only one of the following:
1. $`G_0=1`$;
2. $`G_1=<(1234)>`$;
3. $`G_2=<(13)(24)>`$;
4. $`G_3=<(1432)>`$;
5. $`G_4=<(12)(34)>`$;
6. $`G_5=<(24)>`$;
7. $`G_6=<(14)(23)>`$;
8. $`G_7=<(13)>`$;
call $`L_0,\mathrm{},L_7`$ the corresponding bundles; remark the first 4 are orientable, the last 4 are not.
Now remark that the groups with even index act also on the figure 8 obtained by connecting with arcs edges 1 and 4 and edges 2 and 3; this implies that in the fiber bundles with even index it is possible to substitute the fiber, and so we obtain respectively immersions of a torus in a solid torus, of a Klein bottle in a solid torus, of a torus in a solid Klein bottle and of a Klein bottle in a solid Klein bottle.
Go back to the fiber bundle on the curve of double points. If it is isomorphic to $`L_2`$ or to $`L_4`$ we immerse along a curve homotopically trivial in $`M`$ the fiber bundle in figure 8’s obtained from $`L_2`$ by substitution of fiber and we connect its curve of double points with $`K`$: the curve of double points of the new immersion has tubolar neighborhood isomorphic to $`L_0`$ or to $`L_6`$, respectively, and still represents $`Y`$ up to elements of $`C(M)`$. We can then perform a Rohlin surgery as in figure 8,
and obtain from one side an embedding $`f`$ and from the other a map $`h`$ obtained from the inclusion of the boundary of a tubular neighborhood of a curve representing $`\delta _Y`$ by adding a kink along a longitude $`l`$; and the sum of the class of $`f`$ and $`h`$ differs from $`Y`$ by an element $`gC(M)`$. So if $`\delta _Y`$ is non-orientable we are done. If $`\delta _Y`$ is orientable instead we have to consider the case that $`q_h(l)=2`$. If it is so consider in a disk of $`M`$ an immersion $`h^{}`$ obtained from the standard by adding a kink along a longitude $`l^{}`$ such that $`q_h^{}(l^{})=2`$; if we connect the curve of double points of our immersion to the curve of double points of $`h^{}`$ in the usual manner we obtain an immersion we call again $`h`$ that satisfies $`q_h(l)=0`$.
Now go to the case when the tubular neighborhood of $`K`$ is isomorphic to a $`L_i`$ with $`i`$ odd. We have to consider another auxiliary immersion: take on $`[1,1]`$ the fibration with fibers as in figure 9,
and complete it to the immersion of a closed surface in a disk of $`M`$ by identifying the two figure 8’s at the two edges (without torsions) and close the hole of the fiber on 1 with a 2-disk $`D`$; the resulting immersion has two curves of double points meeting in a triple point in $`D`$; the first has tubular neighborhood isomorphic to $`L_1`$; by connecting this curve to $`K`$ in a new curve $`K^{}`$ we get an immersion $`f_4`$, still representing $`Y`$ up to elements of $`C(M)`$, with the tubular neighborhood of $`K^{}`$ isomorphic to an $`L_i`$ with $`i`$ even; and repeating the previous construction we can assume $`i`$ equal to 0. The other curve of double points of $`f_4`$, say $`K^{\prime \prime }`$, is contained in $`D`$ and is shown in figure 10.
Perform Rohlin surgery on $`K^{}`$: on $`D`$ the result is shown in figure 11:
to the two curves so created one can again perform Rohlin surgery and obtain elements of $`C(M)`$. So, up to again applying the previous construction when $`\delta _Y`$ is orientable, we have a decomposition satisfying the required properties. $`\mathrm{}`$
### 2.3 Other lemmas.
###### Lemma 2.3.1
Let $`f`$ and $`f^{}`$ be two embeddings of compact closed surfaces $`F`$ and $`F^{}`$, respectively, in the same non-orientable 3-manifold $`M`$. Then $`f_cf^{}`$ if and only if $`H_f=H_f^{}`$.
Proof. The “if” part is invariance of $`H_{}`$ up to cobordism, and has already be proved.
For the “only if” part let $`H`$ be $`H_f=H_f^{}`$. If $`M`$ is compact and without boundary $`f`$ and $`f^{}`$ are essentially the loci of zeros of two transverse sections $`s`$ and $`s^{}`$ of the same line bundle $`L`$, Poincaré dual to $`H`$, and this allows to construct the cobordism this way: consider a smooth homotopy $`s_t`$ between $`s`$ and $`s^{}`$, that be constant for $`t`$ close to 0 and to 1; consider on $`M\times I`$ the line bundle pull-back of $`L`$, and consider $`s_t`$ as a section of this line bundle: $`s_t`$ is then transverse and its locus of zeros is the desired cobordism. If $`M`$ is not compact or closed in order to apply Poincaré duality consider a compact 3-manifold $`N`$ with boundary, contained in $`M`$ and containing in its interior a 3-chain that bounds $`f+f^{}`$; consider the compact closed 3-manifold $`\overline{N}`$ obtained by gluing to $`N`$ a second copy of $`N`$ itself; again $`f`$ and $`f^{}`$ represent the same element $`H`$ in the homology of $`\overline{N}`$, so they are loci of zeros of transverse sections $`s`$ and $`s^{}`$ of the same line bundle $`L`$, Poincaré dual (in $`\overline{N}`$) to $`H`$; we can assume that $`s`$ and $`s^{}`$ coincide outside a compact contained in $`N`$; we construct as before a smooth homotopy $`s_t`$, that we pretend relative to a compact containing the complementary of $`N`$; then the construction runs as before, and gives a cobordism which is in fact contained in $`N\times I`$, hence in $`M\times I`$$`\mathrm{}`$
We now prove (2). We first introduce some notation. We denote $`[f]`$ the class of $`f`$ up to reparametrizations of $`F`$, that is, $`g`$ belongs to $`[f]`$ if there exists a diffeomorphism $`\varphi `$ of $`F`$ such that $`g=f\varphi `$; remark that regular homotopy equivalence relation is well defined up to reparametrization, and that immersions that differ by a reparametrization are cobordant (see \[BS95\], pages 657–8). Finally denote by $`[f]K`$ the reparametrization class of $`ff^1(K)`$, that is well defined. The key of the proof are propositions 10.5 and 10.7 of \[BS95\], that we adapt to a non-orientable situation:
###### Lemma 2.3.2
Let $`f`$ be an immersion of $`F`$ in $`M`$. Let $`K`$ be a closed circle in $`f(F)`$ such that $`q_f(f^1(K))=0`$. Then:
1. if $`K`$ is trivial in $`\pi _1(M)`$ then $`[f]`$ and $`[f]K`$ are regularly homotopic;
2. if $`K`$ is trivial in $`H_1(M,𝐙/2𝐙)`$ then $`[f]`$ and $`[f]K`$ are cobordant.
Proof. 1. Let $`c=f^1(K)`$; $`c`$ is orientable in $`F`$. Let $`\varphi `$ be a twist of Dehn of $`F`$ along $`c`$, we prove that $`fc`$ and $`f\varphi `$ are regularly homotopic. We make direct use of the result of \[HH85\].
First remark that, since $`K`$ is homotopically trivial, $`f\varphi `$ is homotopic to $`f`$, that on its side is homotopic to $`fc`$. We can then consider on $`fc`$ and $`f\varphi `$ the correspondence $`C_f`$, as defined in the proof of theorem 1.2.1; we see that $`C_f(fc)=C_f(f\varphi )`$.
Consider an embedding of the standard torus in $`𝐑^3`$ with the $`Spin`$ structure induced by restriction of the unique $`Spin`$ structure on $`𝐑^3`$ as a tubular neighborhood of $`K`$, in such a way that the $`Spin`$ structure of either of the two preimages of $`K`$ in $`\stackrel{~}{M}`$ is respected. The ipothesis $`q_f(c)=0`$ implies that this embedding can be chosen in such a way that $`f`$ restricted to a tubular neighborhood of $`c`$ can be read as the embedding in the standard torus of a band with no half twists. In the same model then $`f\varphi `$ and $`fc`$ are read as shown in figure 12, and far from $`K`$ they coincide with $`f`$.
Now take a representative $`d`$ of an element $`H_1(F,𝐙/2𝐙)`$, transverse to $`c`$: following again the proof of the main theorem of \[HH85\] one shows that if $`dc=0`$ then $`C_f(f\varphi )(d)=C_f(fc)(d)=0`$, and if $`dc=1`$ then $`C_f(f\varphi )(d)=C_f(fc)(d)=1`$; and this ends the proof.
2. The proof of \[BS95\] can be used, by means of 1.$`\mathrm{}`$
We remark that, since immersions that differ by a reparametrization are cobordant, the second part of the lemma implies:
###### Corollary 2.3.3
If $`K`$ is in the image of an immersion $`f`$ and is trivial in $`H_1(M,𝐙/2𝐙)`$ then $`f`$ is cobordant to $`ff^1(K)`$$`\mathrm{}`$
We are now ready for the lemmas that prove (2).
###### Lemma 2.3.4
Let $`\delta H_1(M,𝐙/2𝐙)`$ such that $`w_1(M)(\delta )=0`$; let $`K`$ and $`K^{}`$ be closed circles both representing $`\delta `$; let $`i`$ and $`i^{}`$ be the inclusions of the boundaries of tubular neighborhoods of $`K`$ and $`K^{}`$, respectively. Let $`h`$ and $`h^{}`$ be two immersions obtained from $`i`$ and $`i^{}`$ by adding kinks along longitudes $`l`$ and $`l^{}`$ such that $`q_i(l)=q_i^{}(l^{})=0`$. Then $`h_ch^{}`$.
Proof. This is proven in \[BS95\], page 672. We review their argument. Remark that $`i`$ and $`i^{}`$ are immersions of tori, representing the identity of the semi-group structure; their connected sum $`\overline{i}=i\mathrm{}i^{}`$ still represents the identity. The sum of the cobordism classes of $`h`$ and $`h^{}`$ is represented by their connected sum $`\overline{h}=h\mathrm{}h^{}`$; remark that $`\overline{h}`$ is obtained from $`\overline{i}`$ by adding a kink along a curve $`\overline{l}`$ homologous to $`l+l^{}`$, so that $`q_{\overline{i}}(\overline{l})=q_i(l)+q_i^{}(l^{})=0`$. Moreover the homology class of $`\overline{i}(\overline{l})`$ is $`2\delta `$, that is 0, in $`H_1(M,𝐙/2𝐙)`$. So we can apply corollary 2.3.3 and say that $`\overline{h}`$ is cobordant to $`\overline{i}`$, that is the identity; so $`h^{}`$ belongs to the opposite of the class of $`h`$.
But now do the same construction with two copies of $`h`$, and conclude that $`h`$ itself belongs to its own opposite class, hence the claim. $`\mathrm{}`$
To settle the non-orientable case we first need a remark:
###### Lemma 2.3.5
Let $`K`$ be a curve non-orientable in $`M`$, let $`i`$ be the inclusion of the boundary of a tubular neighborhood $`N`$ of $`K`$, let $`l`$ and $`m`$ be respectively a longitude and the meridian of $`N`$; then
$$il_ci(l+m).$$
Proof. The homotopy class of $`i`$ is odd: in fact $`i`$ can be homotopically deformed to the inclusion of $`K`$, this can slide across a surface representing $`w_1(M)`$ and then going back to $`i`$, and these three steps give a odd self-homotopy of $`i`$. But now $`m`$ represents the dual to $`w_1(N)`$, so that by proposition 1.2.1 adding a kink along $`m`$ doesn’t change regular homotopy class, hence cobordism class. $`\mathrm{}`$
###### Lemma 2.3.6
Let $`\delta H_1(M,𝐙/2𝐙)`$ such that $`w_1(M)(\delta )0`$; let $`K`$ and $`K^{}`$ be closed circles both representing $`\delta `$; let $`i`$ and $`i^{}`$ be the inclusions of the boundaries of tubular neighborhoods of $`K`$ and $`K^{}`$, respectively. Let $`h`$ and $`h^{}`$ be two immersions obtained from $`i`$ and $`i^{}`$ by adding kinks along longitudes $`l`$ and $`l^{}`$. Then $`h_ch^{}`$.
Proof. The inclusions $`i`$ and $`i^{}`$ are immersions of Klein bottles, both representing the identity of the cobordism semi-group. Their connected sum $`\overline{i}=i\mathrm{}i^{}`$ still represents the identity. The sum of the cobordism classes of $`h`$ and $`h^{}`$ is represented by their connected sum $`h\mathrm{}h^{}`$, which is an immersion regularly homotopic to the one obtained from $`\overline{i}`$ by adding a kink along a curve $`\overline{l}`$ homologous to $`l+l^{}`$; this is cobordant to the immersion obtained from $`\overline{i}`$ by adding a kink along the curve $`\overline{l}^{}`$ homologous to $`l+l^{}+m`$, $`m`$ being the meridian of the tubular neighborhood of $`K`$, because of lemma 2.3.5. But both of $`\overline{l}`$ and $`\overline{l}^{}`$ represent $`2\delta `$, hence 0, in $`H_1(M,𝐙/2𝐙)`$, and one of $`q_{\overline{i}}(\overline{l})`$ or $`q_{\overline{i}}(\overline{l}^{})`$ must be 0, see figure 13.
So we can apply corollary 2.3.3, and obtain that $`\overline{h}`$ is cobordant to $`\overline{i}`$, that is the identity class. This implies that $`h^{}`$ belongs to the opposite class to $`h`$.
But now applying the same construction to two copies of $`h`$ we obtain that $`h`$ itself belongs to its own opposite class, hence the claim. $`\mathrm{}`$
We are left now to the part that is in some way more characteristic of the non-orientable case, that is, the point that gives account of the fact that the cobordism groups of non-orientable 3-manifolds are “smaller” than the groups of the orientable case. The key is again the fact that in a non-orientable environment there is a kind of isotopy that doesn’t exist in an orientable manifold, that is, an isotopy that reverses orientation.
###### Proposition 2.3.7
Let $`C(M)`$ be the subgroup of $`N_2(M)`$ given by classes which admit a representative immersed in a disk of $`M`$. $`C(M)`$ is isomorphic to $`𝐙/2𝐙`$.
Proof. Consider a coordinate chart $`(U,\varphi )`$ of $`M`$. The diffeomorphism $`\varphi `$ induces a homomorphism
$$\begin{array}{cccc}\varphi ^{}:& N_2(𝐑^3)& & C(M)\hfill \\ & f& & \varphi f\hfill \end{array}$$
that is evidently surjective. Recall that $`N_2(𝐑^3)`$ is a cyclic group of order 8 generated by the right immersion of Boy of $`𝐏^2`$; we call $`\gamma `$ this immersion. This implies that $`C(M)`$ is a cyclic group generated by $`\varphi ^{}(\gamma )`$, that can’t be trivial since $`𝐏^2`$ generates the cobordism group of surfaces.
Let now $`f_t`$ be a self-isotopy of $`B`$ that reverses the orientation, for example that crosses once a surface representing $`w_1(M)`$. Recall that the inverse of $`\gamma `$ in $`N_2(𝐑^3)`$ is given by composition with a reflection in a plane, so that $`f_1\varphi ^{}(\gamma )`$ is the image via $`\varphi ^{}`$ of the canonical representative of the opposite class to $`\gamma `$, hence represents the opposite class to $`\varphi ^{}(\gamma )`$. But $`f_1\varphi ^{}(\gamma )`$, being isotopic to $`\varphi ^{}(\gamma )`$, is cobordant to it, so that $`\gamma `$ belongs to its own opposite class, hence has order 2. $`\mathrm{}`$
###### Lemma 2.3.8
Let $`g`$ and $`g^{}`$ be immersions whose image is contained in a disk of $`M`$; then $`g_cg^{}`$ if and only if $`n_g=n_g^{}`$.
Proof. The “if” part is invariance of $`n_{}`$ up to cobordism, and has already be proved. For the “only if” it is enough to remark that $`n_{}`$ realizes the homomorphism of proposition 2.3.7$`\mathrm{}`$
Rosa Gini, gini@dm.unipi.it,
Dipartimento di Matematica “Leonida Tonelli”,
via Filippo Buonarroti 2,
I–56127 Pisa, Italy. |
warning/0002/gr-qc0002002.html | ar5iv | text | # Gravitation, the Quantum, and Bohr’s Correspondence Principle
## Abstract
The black hole combines in some sense both the “hydrogen atom” and the “black-body radiation” problems of quantum gravity. This analogy suggests that black-hole quantization may be the key to a quantum theory of gravity. During the last twenty-five years evidence has been mounting that black-hole surface area is indeed quantized, with uniformally spaced area eigenvalues. There is, however, no general agreement on the spacing of the levels. In this essay we use Bohr’s correspondence principle to provide this missing link. We conclude that the fundamental area unit is $`4\mathrm{}\mathrm{ln}3`$. This is the unique spacing consistent both with the area-entropy thermodynamic relation for black holes, with Boltzmann-Einstein formula in statistical physics and with Bohr’s correspondence principle.
Everything in our past experience in physics tells us that general relativity and quantum theory must be approximations, special limits of a single, universal theory. However, despite the flurry of research, which dates back to the 1930s, we still lack a complete theory of quantum gravity. It is believed that black holes may play a major role in our attempts to shed some light on the nature of a quantum theory of gravity (such as the role played by atoms in the early development of quantum mechanics).
The quantization of black holes was proposed long ago in the pioneering work of Bekenstein . The idea was based on the remarkable observation that the horizon area of nonextremal black holes behaves as a classical adiabatic invariant. In the spirit of Ehrenfest principle, any classical adiabatic invariant corresponds to a quantum entity with a discrete spectrum, Bekenstein conjectured that the horizon area of a quantum black hole should have a discrete eigenvalue spectrum.
To elucidate the spacing of the area levels it is instructive to use a semiclassical version of Christodoulou’s reversible processes. Christodoulou showed that the assimilation of a neutral (point) particle by a (nonextremal) black hole is reversible if it is injected at the horizon from a radial turning point of its motion. In this case the black-hole surface area is left unchanged and the changes in the other black-hole parameters (mass, charge, and angular momentum) can be undone by another suitable (reversible) process. (This result was later generalized by Christodoulou and Ruffini for charged point particles ).
However, in a quantum theory the particle cannot be both at the horizon and at a turning point of its motion; this contradicts the Heisenberg quantum uncertainty principle. As a concession to a quantum theory Bekenstein ascribes to the particle a finite effective proper radius $`b`$. This implies that the capture process (of a neutral particle) involves an unavoidable increase $`(\mathrm{\Delta }A)_{\mathrm{min}}`$ in the horizon area :
$$(\mathrm{\Delta }A)_{\mathrm{min}}=8\pi (\mu ^2+P^2)^{1/2}b,$$
(1)
where $`\mu `$ and $`P`$ are the rest mass and physical radial momentum (in an orthonormal tetrad) of the particle, respectively. In the classical case the limit $`b0`$ recovers Christodoulou’s result $`(\mathrm{\Delta }A)_{\mathrm{min}}=0`$ for a reversible process. However, a quantum particle is subjected to a quantum uncertainty – the particle’s center of mass cannot be placed at the horizon with accuracy better than the radial position uncertainty $`\mathrm{}/(2\delta P)`$. This yields a lower bound on the increase in the black-hole surface area due to the assimilation of a (neutral) test particle
$$(\mathrm{\Delta }A)_{\mathrm{min}}=4\pi l_{p}^{}{}_{}{}^{2},$$
(2)
where $`l_p=\left(\frac{G}{c^3}\right)^{1/2}\mathrm{}^{1/2}`$ is the Planck length (we use gravitational units in which $`G=c=1`$). Thus, for nonextremal black holes there is a universal (i.e., independent of the black-hole parameters) minimum area increase as soon as one introduces quantum nuances to the problem.
The universal lower bound Eq. (2) derived by Bekenstein is valid only for neutral particles . Expression (1) can be generalized for a charged particle of rest mass $`\mu `$ and charge $`e`$. Here we obtain
$$(\mathrm{\Delta }A)_{min}=\{\begin{array}{cc}4\pi [2(\mu ^2+P^2)^{1/2}be\mathrm{\Xi }_+b^2],\hfill & b<b^{},\hfill \\ 4\pi (\mu ^2+P^2)/e\mathrm{\Xi }_+,\hfill & bb^{},\hfill \end{array}$$
(3)
where $`\mathrm{\Xi }_+`$ is the black-hole electric field (we assume that $`e\mathrm{\Xi }_+>0`$) and $`b^{}(\mu ^2+P^2)^{1/2}/e\mathrm{\Xi }_+`$. Evidently, the increase in black-hole surface area can be minimized by maximizing the black-hole electric field. Is there a physical mechanism which can prevents us from making expression (3) as small as we wish ? The answer is “yes” ! Vacuum polarization effects set an upper bound to the strength of the black-hole electric field; the critical electric field $`\mathrm{\Xi }_c`$ for pair-production of particles with rest mass $`\mu `$ and charge $`e`$ is $`\mathrm{\Xi }_c=\pi \mu ^2/e\mathrm{}`$ . Therefore, the minimal black-hole area increase is given by
$$(\mathrm{\Delta }A)_{\mathrm{min}}=4l_{p}^{}{}_{}{}^{2}.$$
(4)
Remarkably, this lower bound is independent of the black-hole parameters.
The underling physics which excludes a completely reversible process (for neutral particles) is the Heisenberg quantum uncertainty principle . However, for charged particles it must be supplemented by another physical mechanism – a Schwinger discharge of the black hole (vacuum polarization effects). Without this physical process one could have reached the reversible limit. It seems that nature has “conspired” to prevent this.
It is remarkable that the lower bound found for charged particles is of the same order of magnitude as the one given by Bekenstein for neutral particles, even though they emerge from different physical mechanisms. The universality of the fundamental lower bound (i.e., its independence on the black-hole parameters) is clearly a strong evidence in favor of a uniformly spaced area spectrum for quantum black holes. Hence, one concludes that the quantization condition of the black-hole surface area should be of the form
$$A_n=\gamma l_{p}^{}{}_{}{}^{2}n;n=1,2,\mathrm{},$$
(5)
where $`\gamma `$ is a dimensionless constant.
It should be recognized that the precise values of the universal lower bounds Eqs. (2) and (4) can be challenged. This is a direct consequence of the inherent fuzziness of the uncertainty relation. Nevertheless, it should be clear that the fundamental lower bound must be of the same order of magnitude as the one given by Eq. (4); i.e., we must have $`\gamma =O(4)`$. The small uncertainty in the value of $`\gamma `$ is the price we must pay for not giving our problem a full quantum treatment. In fact, the above analyses are analogous to the well known semiclassical derivation of a lower bound to the ground state energy of the hydrogen atom (calculated by using Heisenberg’s uncertainty principle, without solving explicitly the Schrödinger wave equation). The analogy with usual quantum physics suggests the next step – a wave analysis of black-hole perturbations.
The evolution of small perturbations of a black hole are governed by a one-dimensional Schrödinger-like wave equation (assuming a time dependence of the form $`e^{iwt}`$) . Furthermore, it was noted that, at late times, all perturbations are radiated away in a manner reminiscent of the last pure dying tones of a ringing bell . To describe these free oscillations of the black hole the notion of quasinormal modes was introduced . The quasinormal mode frequencies (ringing frequencies) are characteristic of the black hole itself.
It turns out that there exist an infinite number of quasinormal modes for $`n=0,1,2,\mathrm{}`$ characterizing oscillations with decreasing relaxation times (increasing imaginary part) . On the other hand, the real part of the frequency approaches a constant value as $`n`$ is increased.
Our analysis is based on Bohr’s correspondence principle (1923): “transition frequencies at large quantum numbers should equal classical oscillation frequencies”. Hence, we are interested in the asymptotic behavior (i.e., the $`n\mathrm{}`$ limit) of the ringing frequencies. These are the highly damped black-hole oscillations frequencies, which are compatible with the statement (see, for example, ) “quantum transitions do not take time” (let $`w=w_Riw_I`$, then $`\tau w_{I}^{}{}_{}{}^{1}`$ is the effective relaxation time for the black hole to return to a quiescent state. Hence, the relaxation time $`\tau `$ is arbitrarily small as $`n\mathrm{}`$).
Nollert found that the asymptotic behavior of the ringing frequencies of a Schwarzschild black hole is given by
$$Mw_n=0.0437123\frac{i}{4}\left(n+\frac{1}{2}\right)+O\left[(n+1)^{1/2}\right].$$
(6)
It is important to note that the asymptotic limit is independent of the multipole index $`l`$ of the perturbation field. This is a crucial feature, which is consistent with the interpretation of the highly damped ringing frequencies (in the $`n1`$ limit) as being characteristics of the black hole itself. The asymptotic behavior Eq. (6) was later verified by Andersson using an independent analysis.
We note that the numerical limit $`Re(Mw_n)0.0437123`$ (as $`n\mathrm{}`$) agrees (to the available data given in ) with the expression $`\mathrm{ln}3/(8\pi )`$. This identification is supported by thermodynamic and statistical physics arguments discussed below. Using the relations $`A=16\pi M^2`$ and $`dM=E=\mathrm{}w`$ one finds $`\mathrm{\Delta }A=4l_{p}^{}{}_{}{}^{2}\mathrm{ln}3`$. Thus, we conclude that the dimensionless constant $`\gamma `$ appearing in Eq. (5) is $`\gamma =4\mathrm{ln}3`$ and the area spectrum for a quantum black hole is given by
$$A_n=4l_{p}^{}{}_{}{}^{2}\mathrm{ln}3n;n=1,2,\mathrm{}.$$
(7)
This result is remarkable from a statistical physics point of view ! The semiclassical versions of Christodoulou’s reversible processes, which naturally lead to the conjectured area spectrum Eq. (5), are at the level of mechanics, not statistical physics. In other words, these arguments did not relay in any way on the well known thermodynamic relation between black-hole surface area and entropy. In the spirit of Boltzmann-Einstein formula in statistical physics, Mukhanov and Bekenstein relate $`g_nexp[S_{BH}(n)]`$ to the number of microstates of the black hole that correspond to a particular external macrostate ($`S_{BH}`$ being the black-hole entropy). Namely, $`g_n`$ is the degeneracy of the $`n`$th area eigenvalue. The accepted thermodynamic relation between black-hole surface area and entropy can be met with the requirement that $`g_n`$ has to be an integer for every $`n`$ only when
$$\gamma =4\mathrm{ln}k;k=2,3,\mathrm{}.$$
(8)
Thus, statistical physics arguments force the dimensionless constant $`\gamma `$ in Eq. (5) to be of the form Eq. (8). Still, a specific value of $`k`$ requires further input, which was not aveliable so far. The correspondence principle provides a first independent derivation of the value of $`k`$. It should be mentioned that following the pioneering work of Bekenstein a number of independent calculations (most of them in the last few years) have recovered the uniformally spaced area spectrum Eq. (5) . However, there is no general agreement on the spacing of the levels. Moreover, non of these calculations is compatible with the relation $`\gamma =4\mathrm{ln}k`$, which is a direct consequence of the accepted thermodynamic relation between black-hole surface area and entropy.
The fundamental area spacing $`4l_{p}^{}{}_{}{}^{2}\mathrm{ln}3`$ is the unique value consistent both with the area-entropy thermodynamic relation, with statistical physics arguments (namely, with the Boltzmann-Einstein formula), and with Bohr’s correspondence principle.
ACKNOWLEDGMENTS
It is a pleasure to thank Jacob D. Bekenstein and Tsvi Piran for stimulating discussions. This research was supported by a grant from the Israel Science Foundation. |
warning/0002/hep-th0002080.html | ar5iv | text | # Factorization of correlation functions in coset conformal field theories.
(February,2000)
## Abstract
We use the conformal Ward identities to study the structure of correlation functions in coset conformal field theories. For a large class of primary fields of arbitrary g/h theory a factorization anzatz is found.Corresponding correlation functions are explicitly expressed in terms of correlation functions of two independent WZNW theories for g and h.
Coset theories is an important subclass of two-dimensional conformal invariant QFT’s.(For a review see e.g..) The $`g/h`$ coset theory is based on the Virasoro algebra generated by
$`K(m)=L^g(m)L^h(m),mZ.`$ (1)
The operator $`L^g(m)`$ is a conformal generator of the Wess-Zumino-Novikov-Witten (WZNW) theory \[4-6\] for the Lie algebra $`g`$ and $`hg.`$ In this paper we study the connection between the coset and WZNW theories which follows from (1) at the level of correlation functions and primary fields.
In the case of $`su(2)/u(1)`$ correlation functions of primary fields may be written in terms of correlation functions of the independent WZNW theories for $`su(2)`$ and $`u(1)`$ .Correlation functions of the $`g/u(1)^d,d=1\mathrm{}rankg,`$ cosets have a similar structure.In some correlation functions of minimal models were expressed in terms of correlation functions of two WZNW theories.
In this paper we show that a large class of correlation functions of arbitrary $`g/h`$ coset conformal field theory can be expressed in terms of correlation functions of two independent WZNW theories for $`g`$ and $`h.`$ To find correlation functions of coset primary fields we use the conformal Ward identities .We propose an anzatz for coset primary fields and show that the corresponding correlation functions satisfy the Ward identities. Different factorization properties of $`g/h`$ coset correlation functions were found in .
The results of this paper are in agreement with those of refs. where the $`g/h`$ WZNW model was studied by using path integral approach.
We begin with the affine Lie algebra $`\widehat{g}_k`$ for simple $`g`$
$`[J^a(m),J^b(n)]=if^{abc}J^c(m+n)+km\delta {}_{}{}^{ab}\delta _{m+n,0}^{},`$
where $`f^{abc}`$ are the structure constants of $`g`$ and $`k`$ is the central charge.
The conformal generator $`L^g(m)`$ is given by
$`L^g(m)={\displaystyle \frac{1}{2k+Q_g}}{\displaystyle \underset{n}{}}:J^a(mn)J^a(n):,`$
where $`Q_g`$ is the quadratic Casimir in the adjoint representation of $`g`$.These operators satisfy the commutator relations
$`[L^g(m),L^g(n)]=(mn)L^g(m+n)+c^g\left[{\displaystyle \frac{1}{12}}(m^3m)\right]\delta _{m,n},`$ (2)
$`c^g={\displaystyle \frac{2kdimg}{2k+Q_g}},`$
where $`c^g`$ is the central charge.
Let $`G_R(z)`$ be the primary field of $`\widehat{g}_k`$
$`[J^a(m),G_R(z)]=z^mG_R(z)t_R^a,`$ (3)
$`[t_R^a,t_R^b]=if^{abc}t_R^c,`$
where $`t_R^a`$ is the representation of the generators of $`g`$ for the field $`G_R(z).`$ In the WZNW theory $`G_R(z)`$ also is the primary field of the Virasoro algebra (2)
$`[L^g(m),G_R(z)]=z^{m+1}_zG_R(z)+\mathrm{\Delta }_R(m+1)z^mG_J(z),`$
$`\mathrm{\Delta }_R={\displaystyle \frac{Q_J}{2k+Q_g}},`$
where $`Q_R`$ is the quadratic Casimir of $`g`$ in the representation $`R.`$ Here and in what follows we treat only the holomorphic part.
Let $`\widehat{h}_k`$ be a subalgebra of $`\widehat{g}_k.`$ We assume that it is generated by $`J^A(m),A=1\mathrm{}dimh.`$ The field $`G_R`$ can be decomposed in the set of some irreducible representations of $`h`$
$`G_R(z)={\displaystyle \underset{l}{}}G_{Rl}(z)={\displaystyle \underset{l}{}}P_lG_R(z),`$ (4)
where $`G_{Rl}(z)`$ belongs to the $`l^{}`$s representation and $`P_l`$ is the corresponding projector. The field $`G_{Rl}`$ satisfies the equation
$`[J^A(m),G_{Rl}(z)]=z^mG_{Rl}(z)t_l^A,`$ (5)
where $`t_l^A`$ is the representation of the generators of $`h`$ for the field $`G_{Rl}(z)`$ As well as $`G_R(z)`$ the field $`G_{Rl}(z)`$ is the primary field of the Virasoro algebra (2)
$`[L^g(m),G_{Rl}(z)]=z^{m+1}_zG_{Rl}(z)+\mathrm{\Delta }_R(m+1)z^mG_{Rl}(z),`$ (6)
Correlation functions of these fields can be computed using correlation functions of the WZNW theory
$`<G_{R_1l_1}(z_1)\mathrm{}G_{R_Nl_N}(z_N)>={\displaystyle \underset{i=1}{\overset{N}{}}}P_{l_i}<G_{R_1}(z_1)\mathrm{}G_{R_N}(z_N)>`$
The coset conformal generators $`K(m)`$ (1) satisfy commutator relations (2) with the central charge $`c^{g/h}=c^gc^h`$ . We shall need the relation
$`[K(m),G_{Rl}(z)]=z^{m+1}(_zG_{Rl}(z){\displaystyle \frac{2}{2k+Q_h}}:J^A(z)G_{Rl}(z):t_l^A)`$
$`+\mathrm{\Delta }_{Rl}(m+1)z^mG_{Rl}(z),`$ (7)
where
$`:J^A(z)G_{Rl}(z):={\displaystyle \underset{m<0}{}}J^A(m)z^{m1}G_{Rl}(z)+G_{Rl}(z){\displaystyle \underset{m0}{}}J^A(m)z^{m1}.`$
$`\mathrm{\Delta }_{Rl}`$ is given by
$`\mathrm{\Delta }_{Rl}=\mathrm{\Delta }_R{\displaystyle \frac{Q_l}{2k+Q_h}},`$ (8)
where $`Q_l`$ is the quadratic Casimir of $`h`$ in the representation $`l.`$
Correlation functions of the coset primary fields $`\varphi _i`$ satisfy the conformal Ward identity
$`<K(z)\varphi _1(z_1)\mathrm{}\varphi _N(z_N)>={\displaystyle \underset{i=1}{\overset{N}{}}}\left\{{\displaystyle \frac{\mathrm{\Delta }_i}{(zz_i)^2}}+{\displaystyle \frac{1}{zz_i}}{\displaystyle \frac{}{z_i}}\right\}<\varphi _1(z_1)\mathrm{}\varphi _N(z_N)>,`$ (9)
where $`K(z)=_mK(m)z^{m2}`$ and $`\mathrm{\Delta }_1,\mathrm{\Delta }_2,\mathrm{},\mathrm{\Delta }_N`$ are dimensions of $`\varphi _1,\varphi _2,\mathrm{},\varphi _N`$, respectively.
To find a solution of this equation we shall use an auxiliary WZNW theory. Let $`\widehat{h}_k^{}`$ be the auxiliary affine Lie algebra
$$[\chi ^A(m),\chi ^B(n)]=if^{ABC}\chi ^C(m+n)+k^{}m\delta {}_{}{}^{AB}\delta _{m+n,0}^{}.$$
(10)
where $`A,B,C=1\mathrm{}dimh.`$ The value of $`k^{}`$ will be defined later. Let $`\mathrm{\Phi }_l`$ be the primary field of the WZNW theory for $`\widehat{h}_k^{}`$
$`[\chi ^A(m),\mathrm{\Phi }_l(z)]=z^m\mathrm{\Phi }_l(z)t_l^A,`$ (11)
$`{\displaystyle \frac{}{z}}\mathrm{\Phi }_l(z)={\displaystyle \frac{2}{2k^{}+Q_h}}:\chi ^A(z)\mathrm{\Phi }_l(z):t_l^A,`$ (12)
where $`t_l^A=(t_l^A)^T`$ and $`Q_h`$ is the quadratic Casimir in the adjoint representation of $`h.`$ Eq.(12) was introduced in ref. .
We look for a solution of eq.(9) in the following factorized form
$`{\displaystyle \underset{\alpha _1\mathrm{}\alpha _N}{}}<G_{R_1l_1}^{\alpha _1}(z_1)\mathrm{}G_{R_Nl_N}^{\alpha _N}(z_N)><\mathrm{\Phi }_{l_1}^{\alpha _1}(z_1)\mathrm{}\mathrm{\Phi }_{l_N}^{\alpha _N}(z_N)>`$
$`<(G{}_{R_1l_1}{}^{}(z_1),\mathrm{\Phi }{}_{l_1}{}^{}(z_1))\mathrm{}(G{}_{R_Nl_N}{}^{}(z_N),\mathrm{\Phi }{}_{l_N}{}^{}(z_N))>,`$ (13)
where $`(,)`$ is the bilinear form
$`(G_{Rl}(z),\mathrm{\Phi }_l(z))={\displaystyle \underset{\alpha =1}{\overset{diml}{}}}G_{Rl}^\alpha (z)\mathrm{\Phi }_l^\alpha (z).`$
We shall denote
$`\stackrel{~}{G}_{Rl}(z)=(G{}_{Rl}{}^{}(z),\mathrm{\Phi }{}_{l}{}^{}(z)).`$ (14)
It follows from (5) and (11) that $`\stackrel{~}{G}_{Rl}(z)`$ commutes with the operator $`\stackrel{~}{J}^A(m)=J^A(m)+\chi ^A(m)`$
$`[\stackrel{~}{J}^A(m),\stackrel{~}{G}_{Rl}(z)]=0.`$ (15)
The vacuum state $`|0>`$ is the joint state of the $`\widehat{g}_k`$ and $`\widehat{h}_k^{}`$ WZNW theories. We shall use the following properties of $`|0>`$
$`<0|\stackrel{~}{J}^A(m<0)=\stackrel{~}{J}^A(m0)|0>=0,`$ (16)
$`<0|K_{}(z)=K_+(z)|0>=0,`$ (17)
where
$`K_{}(z)={\displaystyle \underset{m<1}{}}K(m)z^{m2},K_+(z)={\displaystyle \underset{m1}{}}K(m)z^{m2}.`$
Let us compute the left-hand side of eq.(9) using correlation functions (13). To simplify presentation we assume that $`|z|>|z_1|>\mathrm{}>|z_N|`$. Using eqs.(17), (7) and (12) we get
$`<K(z)\stackrel{~}{G}_{R_1l_1}(z_1)\mathrm{}\stackrel{~}{G}_{R_Nl_N}(z_N)>=<K_+(z)\stackrel{~}{G}_{R_1l_1}(z_1)\mathrm{}\stackrel{~}{G}_{R_Nl_N}(z_N)>`$
$`=\left\{{\displaystyle \frac{\mathrm{\Delta }_{R_1l_1}}{(zz_1)^2}}+{\displaystyle \frac{1}{zz_1}}{\displaystyle \frac{}{z_1}}\right\}<\stackrel{~}{G}_{R_1l_1}(z_1)\mathrm{}\stackrel{~}{G}_{R_Nl_N}(z_N)>`$
$`+<\stackrel{~}{G}_{R_1l_1}(z_1)K_+(z)\mathrm{}\stackrel{~}{G}_{R_Nl_N}(z_N)>+{\displaystyle \frac{1}{zz_1}}<T_{R_1l_1}(z_1)\mathrm{}\stackrel{~}{G}_{R_Nl_N}(z_N)>,`$ (18)
where
$`T_{R_1l_1}(z)={\displaystyle \frac{1}{2k^{}+Q_h}}(G_{R_1l_1}(z)t^A,:\chi ^A(z)\mathrm{\Phi }_{l_1}(z):){\displaystyle \frac{1}{2k+Q_h}}(:J^A(z)G_{R_1l_1}(z):t^A,\mathrm{\Phi }_{l_1}(z)).`$
At $`k^{}=kQ_h`$ the field $`T_{R_1l_1}(z)`$ can be written in the form
$`T_{R_1l_1}(z)={\displaystyle \frac{1}{2k+Q_h}}:\stackrel{~}{J}^A(z)(G_{R_1l_1}(z)t_{l_1}^A,\mathrm{\Phi }_{l_1}(z)):.`$
Due to eqs. (15) and (16) the last term of eq.(18) vanishes
$`<T_{R_1l_1}(z_1)\mathrm{}\stackrel{~}{G}_{R_Nl_N}(z_N)>=0.`$
Proceeding inductively one can show that the correlation function (13) satisfies the Ward identity (9).
From the arguments presented above it follows that $`\stackrel{~}{G}_{Rl}(z)`$ (14) represents the primary field of the $`g/h`$ coset theory which has the conformal dimension (8).We took the fields $`G_{Rl}(z)`$ from the decomposition (4). However to prove the factorization only eqs. (5) and (6) were essentially used. These equations have other solutions which can be used to construct coset primary fields.
To construct coset currents let us consider the field $`J(z)=(J^i(z)),J^i(z)=_mJ_m^iz^{m1},i=dimh+1\mathrm{}dimg.`$ It can be decomposed in the set of some irreducible representations of $`h`$
$`J(z)={\displaystyle \underset{s}{}}J_s(z).`$
The field $`J_s(z)`$ satisfies eq.(5) for some $`t_s^A`$ and eq.(6) with the conformal dimension $`\mathrm{\Delta }=1.`$ According to eq. (14) the coset current corresponding to $`J_s(z)`$ is given by
$`\stackrel{~}{J}_s(z)=(J_s(z),\mathrm{\Phi }_s(z)).`$ (19)
It follows from (8) that $`\stackrel{~}{J}_s(z)`$ has the conformal dimension
$`1{\displaystyle \frac{Q_s}{2k+Q_h}},`$
where $`Q_s`$ is the quadratic Casimir of $`h`$ in the representation $`s.`$
Let us consider the $`g/u(1)^d,1drankg,`$ coset theory.In this case the primary field $`G_R(z)`$ is decomposed in the set of one-dimensional representations of $`u(1)^d`$
$`G_R(z)={\displaystyle \underset{\mu =1}{\overset{dimR}{}}}G_{R\mu }(z),`$
$`[J^A(m),G_{R\mu }(z)]=\mu ^Az^mG_{R\mu }(z),`$
where $`\mu =(\mu ^A).`$ A solution of eqs.(11),( 12) is given by
$`\mathrm{\Phi }_\mu (z)=:exp({\displaystyle \frac{i}{k^{}}}\mu \phi (z)):,`$ (20)
$$\phi ^A(z)=q^Ai\chi ^A(0)logz+i\underset{n0}{}\frac{\chi ^A(n)}{n}z^n,$$
(21)
where
$$[q^A,\chi ^B(m)]=i\delta ^{AB}\delta _{m,0}.$$
(22)
According to eqs.(14) and (8) at $`k^{}=k`$ $`\stackrel{~}{G}_{R\mu }(z)=G_{R\mu }(z)\mathrm{\Phi }_\mu (z)`$ represents the coset primary field which has the dimension $`\mathrm{\Delta }_{R\mu }=\mathrm{\Delta }_R\mu ^2/2k.`$
The correlation function of these fields is given by
$`<G_{R_1\mu _1}(z_1)\mathrm{}G_{R_N\mu _N}(z_N)>=<G_{R_1}(z_1)\mathrm{}G_{R_N}(z_N)>{\displaystyle \underset{i<j}{\overset{N}{}}}(z_iz_j)^{\frac{\mu _i\mu _j}{k}}.`$
This is in agreement with the results of refs.. Parafermion $`g/u(1)^d`$ currents in the form (19) were obtained in .
The results presented in this paper can be extended in many directions. The most important is to study the factorization properties of the $`W/h`$ coset conformal field theory. It is also interesting to find primary fields and describe the corresponding operator algebra. This is presently being studied. |
warning/0002/nucl-th0002056.html | ar5iv | text | # 𝜂𝑑 scattering in the region of the 𝑆₁₁ resonance
## I Introduction
The elastic scattering $`\eta d`$ reaction has been studied recently by several authors in order to investigate the existence of a resonance or a quasibound state in this system, for which there are experimental indications . Some of these studies concluded that such an state would exist for certain values of the two-body $`\eta `$N data. However, since they used in one form or another incomplete information on the two-body subsystems (in particular that corresponding to the $`\eta N`$ sector), we believe that a new calculation is required which takes into account all the information that is now available. In particular, we study the effect of the repulsion at short distances of the $`NN`$ interaction and take into account the $`\eta N\eta N`$ scattering amplitude that has been determined recently .
We will present the Faddeev formalism in section II. In section III we will first calculate the eta-deuteron scattering length to compare with the results of multiple-scattering theories as well as with separable potentials models and finally we will present the predictions of our model for $`\eta d`$ scattering. We will give our conclusions in section IV.
## II Faddeev Formalism
Let us consider a system of three particles, where two of them are identical, interacting pairwise through separable potentials that act only in S-waves. In the case of the $`\eta d`$ system, S-wave means, for the eta-nucleon pair the $`S_{11}`$ channel, and for the nucleon-nucleon pair the $`{}_{}{}^{3}S_{1}^{}`$ channel. The two-body T-matrix of the pair $`jk`$ will be assumed of the separable form
$$t_i(p_i,p_i^{};E)=g_i(p_i)\tau _i(E)g_i(p_i^{}),$$
(1)
where $`\tau _i(E)`$ and $`g_i(p_i)`$ will be specified later. In the following we will identify particle 1 with the $`\eta `$ and the identical particles 2 and 3 with the two nucleons.
The Faddeev equations for the case of $`\eta d`$ scattering can be solved explicitely for the $`\eta dN(N\eta )`$ transition amplitude $`T_2`$, with nucleon 2 being the spectator particle in the final state. One obtains
$$T_2(q_2;E)=2K_{21}(q_2,q_{10};E)+_0^{\mathrm{}}q_{2}^{}{}_{}{}^{2}𝑑q_2^{}M(q_2,q_2^{};E)\tau _2(Eq_{2}^{}{}_{}{}^{2}/2\nu _2)T_2(q_2^{};E),$$
(2)
where
$$M(q_2,q_2^{};E)=K_{23}(q_2,q_2^{};E)+2_0^{\mathrm{}}q_1^2𝑑q_1K_{21}(q_2,q_1;E)\tau _1(Eq_1^2/2\nu _1)K_{12}(q_1,q_2^{};E),$$
(3)
and the driving term for the transition from the amplitude with spectator j in the initial state to the one with spectator i in the final state is
$$K_{ij}(q_i,q_j;E)=\frac{1}{2}_1^1𝑑cos\theta \frac{g_i(p_i)g_j(p_j)}{Ep_i^2/2\mu _iq_i^2/2\nu _i+iϵ},$$
(4)
with
$$p_i=\left(\frac{\mu _i^2}{m_k^2}q_i^2+q_j^2+2\frac{\mu _i}{m_k}q_iq_jcos\theta \right)^{1/2},$$
(5)
$$p_j=\left(\frac{\mu _j^2}{m_k^2}q_j^2+q_i^2+2\frac{\mu _j}{m_k}q_iq_jcos\theta \right)^{1/2}.$$
(6)
$`\mu _i`$ and $`\nu _i`$ are the reduced masses
$$\mu _i=\frac{m_jm_k}{m_j+m_k},$$
(7)
$$\nu _i=\frac{m_i(m_j+m_k)}{m_i+m_j+m_k},$$
(8)
and the $`\eta d`$ on-shell relative momentum $`q_{10}`$ is defined by the relation
$$E=\frac{q_{10}^2}{2\nu _1}B_d,$$
(9)
where $`B_d`$ is the binding energy of the deuteron.
The $`\eta d`$ elastic-scattering amplitude is obtained from the solution of Eq. (2) by inserting into it a final $`\eta d`$ state, as
$$F_{\eta d}(q_{10})=\frac{\pi \nu _1}{N}_0^{\mathrm{}}q_2^2𝑑q_2K_{12}(q_{10},q_2;E)\tau _2(Eq_2^2/2\nu _2)T_2(q_2;E).$$
(10)
where $`N`$ is the normalization of the deuteron wave function
$$N=_0^{\mathrm{}}p_1^2𝑑p_1\left[\frac{g_1(p_1)}{B_d+p_1^2/2\mu _1}\right]^2.$$
(11)
The $`\eta d`$ scattering length is given by
$$A_{\eta d}=F_{\eta d}(0),$$
(12)
while the integrated elastic cross section is given by
$$\sigma _{ELAS}=4\pi |F_{\eta d}(q_{10})|^2.$$
(13)
Notice that Eqs. (2) - (4) and (10) do not include the $`\pi N`$ channel explicitely, but only through the inelasticity of the $`\eta N`$ channel. As for the the $`\eta N`$ inelasticity due to the $`\pi \pi N`$ channel, its contribution is not yet included at this stage of the calculations.
## III Results
We started by calculating the $`\eta d`$ scattering length (12), by solving the integral equation (2) with the method of matrix inversion after replacing the integration with a gaussian mesh. Notice that in this case Eqs. (2)-(4) are free of singularities ( the singularity of $`\tau _1`$ at $`E=B_d`$ in eq. (3) occurs for $`q_1=0`$ and is regularized by the integration volume element). Additionally, we also calculated the integrated elastic cross section of $`\eta d`$ scattering. In order to solve the integral equation (2) above threshold, we used the method of contour rotation .
### A $`A_{\eta d}`$ with non-dynamical separable models
We will calculate here the $`\eta d`$ scattering length $`A_{\eta d}`$ for the models proposed in Refs. . The signal that a quasibound state exists for a given model is that the real part of $`A_{\eta d}`$ becomes negative while the imaginary part gets large.
Notice that in Eq. (1) we have assumed a separable model for the two-body amplitudes $`t_i`$. This form of the T-matrix is obtained if one assumes a separable potential between particles $`j`$ and $`k`$: a dynamical two-body equation determines the function $`\tau _i(E)`$ by the form factors $`g_i(p_i)`$, which carry information on the the range of the potential, and the strength parameter of the potential. However, in the AGS formalism used in Ref. and in the multiple-scattering approach used in Ref. the function $`\tau _2(E)`$ for the $`\eta N`$ subsystem, instead of being calculated, has been chosen independently of the form factor $`g_2(p_2)`$. Such an assumption violates the spirit of the Faddeev approach which requires a two-body interaction in order to relate the off-shell behavior of the T-matrix in the energy variable $`E`$ to the off-shell behavior in the momentum variable $`p_2`$. Nevertheless, it is instructive to repeat those calculations in order to check the accuracy of our numerical solution by comparing with the exact result of Ref. as well as to test the convergence of the multiple-scattering schemme developed in which is based in a partial summation of the multiple-scattering series.
In both Refs. the form factor $`g_2(p_2)`$ has been taken of the Yamaguchi form
$$g_2(p_2)=\frac{1}{\alpha _2^2+p_2^2},$$
(14)
with $`\alpha _2=3.316`$ fm$`{}_{}{}^{}1`$. In Ref. , the function $`\tau _2(E)`$ has been parametrized as
$$\tau _2(E)=\frac{\lambda _\eta }{EE_0+\frac{i}{2}\mathrm{\Gamma }},$$
(15)
with $`E_0=1535`$ MeV - $`(m_N+m_\eta )`$ and $`\mathrm{\Gamma }=150`$ MeV. The parameter $`\lambda _\eta `$ in Eq. (15) was chosen to reproduce the complex $`\eta N`$ scattering length $`a_{\eta N}`$, by using the relation
$$t_2(0,0;0)=\frac{a_{\eta N}}{\pi \mu _2},$$
(16)
where $`\mu _2`$ is the $`\eta N`$ reduced mass. This leads to
$$\lambda _\eta =\frac{\alpha _2^4(E_0i\mathrm{\Gamma }/2)}{\pi \mu _2}a_{\eta N}.$$
(17)
As for the nucleon-nucleon separable T-matrix used in Ref. , it was generated from a Yamaguchi separable potential with an energy-dependent strength . Using these parameters we calculated the $`\eta d`$ scattering length $`A_{\eta d}`$ with the formalism described in the previous section for a variety of values of the $`\eta N`$ scattering length $`a_{\eta N}`$ that have been proposed in the literature (see Refs. ). We show the results of this comparison in table I. The results of Ref. using the AGS formalism are given in column two and the ones obtained with the formalism of the previous section are given in column three. As one can see from table I there is very good agreement between the two calculations. The small discrepancies shown in some cases are of no significance since they occur for the values of $`a_{\eta N}`$ allowing the quasibound state to occur and thereby the solutions are highly unstable. In our case, in this situation we had to use a large number of mesh points in order to guarantee stability.
In the multiple-scattering approach of Ref. an approximate formula was used which is based in a partial sumation of the multiple-scattering series. The function $`\tau _2(E)`$ was taken to be constant
$$\tau _2(E)=\lambda _\eta .$$
(18)
Using the relation (16) this gives
$$\lambda _\eta =\frac{\alpha _2^4}{\pi \mu _2}a_{\eta N},$$
(19)
which will be referred to as their model I. They used also a second model which will be referred to as model II for which instead of Eq. (19) they took
$$\lambda _\eta =\frac{\alpha _2^4}{\pi \mu _2}\frac{a_{\eta N}}{1iq_0a_{\eta N}},$$
(20)
with $`q_0`$ = i0.367 fm<sup>-1</sup>. For the nucleon-nucleon interaction they used a Yamaguchi separable potential with a range parameter $`\alpha _1=1.41`$ fm$`{}_{}{}^{}1`$.
We compare in table II the results of our exact calculations which we obtained using the parameters of Ref. with their results using an approximate formula for the two models I and II. As it can be seen from this table, the approximate formula of the multiple scattering series works very well for small values of $`a_{\eta N}`$, as expected from convergence arguments. When $`a_{\eta N}`$ is large the multiple scattering series formula deviates more from the exact result, nevertheless, it is still qualitatively correct, since it predicts correctly the quasibound states in all the cases where they exist for both models.
### B $`A_{\eta d}`$ with separable-potential models
In the previous subsection we have seen that the models of Refs. predict a quasibound state if the real part of $`a_{\eta N}`$ is of the order of 0.7-0.8 fm. However, since their $`\eta N`$ T-matrix is not derived from a potential their function $`\tau _2(E)`$ is not constrained by their form factor $`g_2(p_2)`$. We will therefore construct separable potential models of the coupled $`\eta N`$ \- $`\pi N`$ system that reproduce in one case just the complex scattering length $`a_{\eta N}`$ for arbitrary values of the $`\eta N`$ range-parameter $`\alpha _2`$ and in another case the full $`\eta N`$-$`\eta N`$ scattering amplitude around the $`S_{11}`$ resonance.
Similarly, we will consider two different models of the $`NN`$ interaction; a simple Yamaguchi potential that does not have short-range repulsion and a PEST model which has the same half-off-shell behavior as the Paris potential so that it contains short-range repulsion.
If we use in the Lippmann-Schwinger equation of the coupled $`\eta N`$ \- $`\pi N`$ system the separable model
$$<p|V_{\eta \eta }|p^{}>=\lambda _\eta g_2(p)g_2(p^{}),$$
(21)
$$<p|V_{\pi \pi }|p^{}>=\lambda _\pi g_\pi (p)g_\pi (p^{}),$$
(22)
$$<p|V_{\eta \pi }|p^{}>=\pm \sqrt{\lambda _\eta \lambda _\pi }g_2(p)g_\pi (p^{}),$$
(23)
the T-matrices are of the form
$$<p|t_{\eta \eta }(E)|p^{}>=g_2(p)\tau _2(E)g_2(p^{}),$$
(24)
$$<p|t_{\pi \pi }(E)|p^{}>=\frac{\lambda _\pi }{\lambda _\eta }g_\pi (p)\tau _2(E)g_\pi (p^{}),$$
(25)
$$<p|t_{\eta \pi }(E)|p^{}>=\pm \sqrt{\frac{\lambda _\pi }{\lambda _\eta }}g_2(p)\tau _2(E)g_\pi (p^{}),$$
(26)
with
$$\frac{1}{\tau _2(E)}=\frac{1}{\lambda _\eta }G_2(E)\frac{\lambda _\pi }{\lambda _\eta }G_\pi (E),$$
(27)
$$G_2(E)=_0^{\mathrm{}}p^2𝑑p\frac{g_2^2(p)}{Ep^2/2\mu _2+iϵ},$$
(28)
$$G_\pi (E)=_0^{\mathrm{}}p^2𝑑p\frac{g_\pi ^2(p)}{E+p_0^2/2\mu _\pi p^2/2\mu _\pi +iϵ}.$$
(29)
$`\mu _2`$ and $`\mu _\pi `$ are the $`\eta N`$ and $`\pi N`$ reduced masses respectively while $`p_0`$ is the $`\pi N`$ relative momentum at the $`\eta N`$ threshold, i.e.,
$$p_0^2=\frac{[s_0(m_\pi +m_N)^2][s_0(m_\pi m_N)^2]}{4s_0},$$
(30)
with
$$s_0=(m_\eta +m_N)^2.$$
(31)
If we use simple Yamaguchi form factors
$$g_2(p)=\frac{1}{\alpha _2^2+p^2},$$
(32)
$$g_\pi (p)=\frac{1}{\alpha _\pi ^2+p^2},$$
(33)
we find that the strengths $`\lambda _\eta `$ and $`\lambda _\pi `$ can be obtained in terms of the real and imaginary parts of $`a_{\eta N}`$ as
$$\lambda _\eta ^1=\frac{\pi \mu _2}{2\alpha _2^3}\frac{\alpha _\pi ^2p_0^2}{2\alpha _\pi p_0}\frac{\pi \mu _2}{\alpha _2^4}\frac{Ima_{\eta N}}{|a_{\eta N}|^2}\frac{\pi \mu _2}{\alpha _2^4}\frac{Rea_{\eta N}}{|a_{\eta N}|^2},$$
(34)
$$\lambda _\pi =\lambda _\eta \frac{\mu _2}{\mu _\pi }\frac{(\alpha _\pi ^2+p_0^2)^2}{p_0\alpha _2^4}\frac{Ima_{\eta N}}{|a_{\eta N}|^2}.$$
(35)
Since we do not include the pion channel explicitly but only through the function $`\tau _2(E)`$ (see Eq. (27)), we will fix the range of the $`\pi N`$ potential to the value $`\alpha _\pi =p_0`$, for which case the second term in the r.h.s. of Eq. (34) drops out and it becomes clear that in this case the strength of the $`\eta N`$ potential is determined by the real part of the $`\eta N`$ scattering length and the strength of the $`\pi N`$ potential is determined by the imaginary part of the $`\eta N`$ scattering length for a given value of the range $`\alpha _2`$. Since this models are based in a Yamaguchi form factor for the $`\eta N`$ potential we will refer to them as $`Y_{\eta N}`$ models.
We constructed also more realistic separable models that reproduce not only the $`\eta N`$ scattering length but also the most important features of the $`S_{11}`$ resonance such as its position and width. For this we considered the $`S_{11}`$ amplitudes obtained from the analyses of Refs. . We found that with a simple Yamaguchi model of the $`\eta N`$ form factor is not possible to generate a resonance in the $`\eta N`$ $`S_{11}`$ channel. We therefore changed the form factor $`g_2(p)`$ instead of Eq. (32) to
$$g_2(p)=\frac{A+p^2}{(\alpha _2^2+p^2)^2},$$
(36)
while keeping for the $`\pi N`$ form factor the Yamaguchi form (33). We give in table III the parameters $`\alpha _\pi `$, $`\lambda _\pi `$, $`\alpha _2`$, $`A`$, $`\lambda _\eta `$ of the coupled $`\eta N`$-$`\pi N`$ separable potentials fitted to the $`S_{11}`$ amplitudes of as well as to the models A, B, C, and D of . We show in Fig. 1, as an example, the $`\eta N`$-$`\eta N`$ amplitude of Ref. (dashed lines) compared with the ones of our separable-potential model. Similar results are obtained for the other models. Since these models generate a resonance in the $`\eta N`$ channel we will refer to them as $`R_{\eta N}`$ models. Since we do not include the pion channel explicitly, only the $`t_{\eta \eta }`$ component of the coupled $`\eta N`$-$`\pi N`$ T-matrix given by Eq. (24) has been used after identifying $`t_{\eta \eta }`$ with $`t_2`$ of Eq. (1).
In the case of the $`NN`$ interaction we have considered two models; the simple Yamaguchi model used in Ref. which has a range parameter $`\alpha _1=1.41`$ fm<sup>-1</sup> (which we will refer to as the $`Y_{NN}`$ model) and the PEST potential constructed in Ref. (which we will refer to as the $`P_{NN}`$ model) that is of the form
$$g_1(p_1)=\underset{n=1}{\overset{6}{}}\frac{C_n}{p_1^2+\beta _n},$$
(37)
where the parameters $`C_n`$ and $`\beta _n`$ are given in Ref. . The half-off-shell T-matrix of this separable potential has the same behavior as that of the Paris potential and therefore it takes into account the repulsion at short distances that is present in the nucleon-nucleon force.
We give in table IV the results of our separable-potential models for the $`\eta d`$ scattering length $`A_{\eta d}`$ where we have considered all four combinations of the $`\eta N`$ and $`NN`$ separable-potential models. In the case of the $`\eta N`$ Yamaguchi model $`Y_{\eta N}`$ we took the range parameter $`\alpha _2=3.316`$ fm<sup>-1</sup> which is the same as in Refs. . Here however the corresponding T-matrix is calculated through the Lippmann-Schwinger equation. In the first column we give the reference for the $`\eta N`$ $`S_{11}`$ amplitude that we used to fit the $`R_{\eta N}`$ model and in the second column we give the $`\eta N`$ scattering length of that amplitude which has been used to construct the $`Y_{\eta N}`$ model. The third column gives $`A_{\eta d}`$ using simple Yamaguchi models for the $`\eta N`$ and $`NN`$ interactions and it shows that even with these simple separable models the quasibound state only appears when $`Rea_{\eta N}`$ is about 1.05-1.07 fm while in the multiple-scattering approaches of the previous section it appeared already with $`Rea_{\eta N}`$ 0.6-0.9 fm. The fourth column gives the results of the Yamaguchi model for the $`\eta N`$ amplitude and the PEST model for the $`NN`$ amplitude and it shows that the $`NN`$ short-range repulsion also works against quasibinding since it wipes out the quasibound state although the $`ImA_{\eta d}`$ remains large. The fifth and sixth columns contain the results of the resonant model of the $`\eta N`$ amplitude with Yamaguchi and PEST models for the $`NN`$ interaction respectively, and they show that the attraction of the system is greatly reduced when one takes into account the resonant nature of the $`\eta N`$ amplitude. Notice, however, that experimentally the effects will be similar whether there is a quasibound state or not since in both cases there will be an enhancement of the cross section at threshold.
### C $`\eta d`$ scattering
We calculated the integrated elastic cross section of $`\eta d`$ scattering in the region of the $`S_{11}`$ resonance for the six resonant models of the $`\eta N`$ interaction given in table III and the realistic PEST potential for the $`NN`$ interaction. We show in Fig. 2 the results of the three-body model (solid lines) and of the impulse approximation (dashed lines). At threshold, the results of all the three-body models are about one order of magnitude larger than those of the impulse approximation while at higher energies they are 2 or 3 times smaller. The behavior of the cross section at threshold indicates that even though the quasibound state is not present, the interaction in this region is very strong since it enhances the impulse approximation result by about one order of magnitude. Thus, a signal of this behavior may appear also in other processes where there is an $`\eta N`$ final state like the $`np\eta d`$ reaction where a large enhancement in the cross section has been observed in the region near threshold .
In order to illustrate the effect of the strong $`\eta d`$ interaction in the reaction $`np\eta d`$ we will estimate the enhancement of the $`np\eta d`$ cross section at threshold due to the $`\eta d`$ rescattering. We write the amplitude of the process $`np\eta d`$ as
$$A=B+BG_0T_{\eta d},$$
(38)
where $`B`$ is the amplitude of the $`np\eta d`$ process without $`\eta d`$ rescattering, $`G_0`$ is the two-body Lippmann-Schwinger propagator of the intermediate $`\eta d`$ state, and $`T_{\eta d}`$ is the half-off-shell T-matrix of the elastic $`\eta d`$ process. If we introduce the $`\eta d`$ elastic-scattering amplitude $`F_{\eta d}=\pi \mu _{\eta d}T_{\eta d}`$, where $`\mu _{\eta d}`$ is the $`\eta d`$ reduced mass, then at threshold, Eq. (38) is written explicitly as
$$A=B(0)[1+\frac{2}{\pi }_0^{\mathrm{}}𝑑q_1\frac{B(q_1)}{B(0)}F_{\eta d}(q_1)],$$
(39)
where $`F_{\eta d}(q_1)`$ is given by Eq. (10) with $`q_{10}`$ replaced by $`q_1`$. Therefore, the enhancement factor of the $`np\eta d`$ cross section due to $`\eta d`$ rescattering is
$$f=|1+\frac{2}{\pi }_0^{\mathrm{}}𝑑q_1\frac{B(q_1)}{B(0)}F_{\eta d}(q_1)|^2.$$
(40)
The amplitude $`B`$ of the $`np\eta d`$ process without $`\eta d`$ rescattering is presumably given by meson exchanges such as $`\pi `$, $`\rho `$, and $`\eta `$ followed by the excitation and decay of the $`S_{11}`$ resonance . The explicit form of the production operator without $`\eta d`$ rescattering is not so important since in our estimate of the enhancement factor given by Eq. (40) only the ratio $`B(q_1)/B(0)`$ enters. Therefore, we take for it the $`\eta `$-exchange amplitude generated by our three-body model, i.e.,
$$B(q_1)=_0^{\mathrm{}}q_2^2𝑑q_2K_{12}(q_1,q_2;E)\tau _2(Eq_2^2/2\nu _2)K_{23}(q_2,q_3;E),$$
(41)
where $`E=B_d`$ and $`q_3=\sqrt{(m_\eta +m_d)^2m_N^2}`$ is the momentum corresponding to an initial $`NN`$ state. Using the models 1-6 of Fig. 2, taken from references and included in Table IV, we obtained the enhancement factors $`f`$ = 2.5, 2.7, 3.1, 3.3, 4.7, and 5.1 respectively. These estimates are quite comparable to the enhancement factors observed in Ref. , in special for the 4 first models.
## IV Conclusions
Recently, the experimental band for the $`\eta N`$ scattering length $`a_{\eta N}`$ has been pushed towards larger values for its real part . For those larger values, the $`\eta N`$ models not generated directly from an integral equation, which would fix naturally their off-energy-shell behavior needed in three-body calculations, predict a quasi-bound $`\eta NN`$ state. If instead, the new data is used to generate $`\eta `$N t-matrices calculated from a potential and an integral equation, our results indicate that a realistic NN interaction, like the Paris potential, through its short-range correlations, prevents the existence of the bound-state, independently of the $`\eta `$N models, provided they have been built dynamically. We confirmed then that the predictions of a $`\eta NN`$ eta-mesic nucleus are crucially affected by the off-shell behavior of the underlying $`\eta `$-N models. Importantly, however, is that even for an inexisting quasi-bound state, an exact three-body calculation for the multi-scattering series in the final state predicts a severe enhancement of the elastic $`\eta d`$ cross-section in the narrow region from threshold to 5-10 MeV above threshold. This result is independent of the $`\eta N`$ two-body models. Very likely, the enhancement predicted by the exact three body calculations is related to the one observed in the reactions $`np\eta d`$ and $`\gamma d\eta d`$ . We actually made an estimate of the enhancement of the cross section of $`np\eta d`$ due to the $`\eta d`$ final state interaction. Within the three-body model of our work, this enhancement is in the ball park of the empirical findings of .
###### Acknowledgements.
This work was supported in part by COFAA-IPN (México) and by Fundação para a Ciência e a Tecnologia, MCT (Portugal) under contracts PRAXIS XXI/BCC/18975/98 and PRAXIS/P/FIS/10031/1998. |
warning/0002/cond-mat0002075.html | ar5iv | text | # Finite-time singularity in the dynamics of the world population, economic and financial indices. Running title: Finite-time singularity in world population growth
## 1 Introduction
Both the world economy as well as the human population have grown at a tremendous pace especially during the last two centuries. It is estimated that 2000 years ago the population of the world was approximately 300 million and for a long time the world population did not grow significantly, since periods of growth were followed by periods of decline. It took more than 1600 years for the world population to double to 600 million and since then the growth has accelerated. It reached 1 billion in 1804 (204 years later), 2 billion in 1927 (123 years later), 3 billion in 1960 (33 years later), 4 billion in 1974 (14 years later), 5 billion in 1987 (13 years later) and 6 billion in 1999 (12 years later). This rapidly accelerating growth has raised sincere worries about its sustainability as well as concerns that we humans as a result might cause severe and irreversible damage to eco-systems, global weather systems etc . At, what one may say the other extreme, the optimists expect that the innovative spirit of mankind will be able to solve the problems associated with a continuing increase in the growth rate . Specifically, they believe that the world economic development will continue as a successive unfolding of revolutions, e.g., the Internet, bio-technological and other yet unknown innovations, replacing the prior agricultural, industrial, medical and information revolutions of the past. Irrespective of the interpretation, the important point is the presence of an acceleration in the growth rate. Here, it is first shown that, contrary to common belief, both the Earth human population as well as its economic output have grown faster that exponential for most of the known history and most strikingly so in the last centuries. Furthermore, we will show that both the population growth rate and the economic growth rate are consistent with a spontaneous singularity at the same critical time $`2052\pm 10`$ and with the same characteristic self-similar geometric patterns (defined below as log-periodic oscillations). Multivariate dynamical equations coupling population, capital and R&D and technology can indeed produce such an “explosion” in the population even though the isolated dynamics do not. In particular, this interplay provides an explanation of our finding of the same value of the critical time $`t_c(2052\pm 10)`$ both for the population and the economic indices. As a consequence, even the optimistic view has to be revised, since the acceleration of the growth rate contains endogenously its own limit in the shape of a finite-time singularity to be interpreted as a transition to a qualitatively new behaviour. Close to the mathematical singularity, finite-size effects will smoothen the transition and it is quite possible that Mankind may already have entered this transition phase. Possible scenarios for the cross-over and the new regime are discussed.
### 1.1 The logistic equation and finite-time singularities
As a standard model of population growth, Malthus’ model assumes that the size of a population increases by a fixed proportion $`\tau `$ over a given period of time independently of the size of the population and thus gives an exponential growth. The logistic equation attempts to correct for the resulting unbounded exponential growth by assuming a finite carrying capacity $`K`$ such that the population instead evolves according to
$$\frac{dp}{dt}=rp(t)\left[Kp(t)\right].$$
(1)
Cohen and others (see and references therein) have put forward idealised models taking into account interaction between the human population $`p(t)`$ and the corresponding carrying capacity $`K(t)`$ by assuming that $`K(t)`$ increases with $`p(t)`$ due to technological progress such as the use of tools and fire, the development of agriculture, the use of fossil fuels, fertilisers etc. as well an expansion into new habitats and the removal of limiting factors by the development of vaccines, pesticides, antibiotics, etc. If $`K(t)>p(t)`$, then $`p(t)`$ explodes to infinity after a finite time creating a singularity. In this case, the limiting factor $`p(t)`$ can be dropped out and, assuming a simple power law relationship $`Kp^\delta `$ with $`\delta >1`$, (1) becomes
$$\frac{dp}{dt}=r[p(t)]^{1+\delta },$$
(2)
where the growth rate accelerates with time according to $`r[p(t)]^\delta `$. The generic consequence of a power law acceleration in the growth rate is the appearance of singularities in finite time:
$$p(t)(t_ct)^z,\mathrm{with}z=\frac{1}{\delta }\text{ and }t\text{ close to }t_c.$$
(3)
Equation (2) is said to have a “spontaneous” or “movable” singularity at the critical time $`t_c`$ , the critical time $`t_c`$ being determined by the constant of integration, i.e., the initial condition $`p(t=0)`$. One can get an intuitive understanding of such singularities by looking at the function $`p(t)=\mathrm{exp}\left(tp(t)\right)`$ which corresponds to replacing $`\tau K`$ by $`p`$ in Malthus’ exponential solution $`p(t)=p(0)\mathrm{exp}[\tau Kt]`$. $`p`$ is then the solution of $`dp/dt=p^2/(1tp)`$ leading to an ever increasing growth with the explicit solution
$$p(t)=e\left(1C\sqrt{t_ct}\right),$$
(4)
where $`t_c=1/e=0.368`$, $`C`$ is a numerical factor and the exponent $`z=1/2`$. In this case, the finite time spontaneous singularity does not lead to a divergence of the population at the critical time $`t_c`$; only the growth rate diverges at $`(t_ct)^{1/2}`$. Spontaneous singularities in ODE’s and PDE’s are quite common and have been found in many well-established models of natural systems either at special points in space such as in the Euler equations of inviscid fluids or in the equations of General Relativity coupled to a mass field leading to the formation of black holes , in models of micro-organisms aggregating to form fruiting bodies , or to the more prosaic rotating coin (Euler’s disk) , see for a review. Some of the most prominent, as well as more controversial, examples due to their impact on human society are models of rupture and material failure , earthquakes and stock market crashes .
### 1.2 Data sets and methodology
Here, we examine several data sets expressing the development of mankind on Earth in term of size and economic impact, to test the hypothesis that our history might be compatible with a future finite-time singularity. These data sets are as follows.
* The human population data from 0 to 1998 was retrieved from the web-site of The United Nations Population Division Department of Economic and Social Affairs (http://www.popin.org/pop1998/).
* The GDP of the world from 0 to 1998, estimated by J. Bradford DeLong at the Department of Economics, U.C. Berkeley , was given to us by R. Hanson .
* The financial data series include the Dow Jones index from 1790 to 2000, the Standard & Poor (S&P) index from 1871 to 2000, as well as a number of regional and global indices since 1920. The Dow Jones index was constructed by The Foundation for the Study of Cycles . It is the Dow Jones index back to 1896, which has been extrapolated back to 1790 and further. The other indices are from Global Financial Data . These indices are constructed as follows. For the S&P, the data from 1871 to 1918 are from the Cowles commission, which back-calculated the data using the Commercial and Financial Chronicle. From 1918, the data is the Standard and Poor’s Composite index (S&P) of stocks. The other indices uses Global Financial Data’s indices from 1919 through 1969 and Morgan Stanley Capital International’s indices from 1970 through 2000. The EAFE Index includes Europe, Australia and the Far East. The Latin America Index includes Argentina, Brazil, Chile, Colombia, Mexico, Peru and Venezuela.
Demographers usually construct population projections in a disaggregated manner, filtering the data by age, stage of development, region, etc. Disaggregating and controlling for such variables are thought to be crucial for demographic development and for any reliable population prediction. Here, we propose a different strategy based on aggregated data, which is justified by the following concept: in order to get a meaningful prediction at an aggregate level, it is often more relevant to study aggregate variables than “local” variables that can miss the whole picture in favor of special idiosyncrasies. To take an example from material sciences, the prediction of the failure of heterogeneous materials subjected to stress can be performed according to two methodologies. Material scientists often analyse in exquisite details the wave forms of the acoustic emissions or other signatures of damage resulting from micro-cracking within the material. However, this is of very little help to predict the overall failure which is often a cooperative global phenomenon resulting from the interactions and interplay between the many different micro-cracks nucleating, growing and fusing within the materials. In this example, it has been shown indeed that aggregating all the acoustic emissions in a single aggregated variable is much better for prediction purpose .
### 1.3 Content of the paper
In the next section, we first show that the exponential model is utterly inadequate in describing the population growth as well as the growth in the World GDP and the global and regional financial indices. We then present the alternative model consisting of a power law growth ending at a critical time $`t_c`$. We first give a non-parametric approach complemented by a fitting procedure. Section 3 proposes a first generalization of power laws with complex exponents, leading to so-called log-periodic oscillations decorating the overall power law acceleration. The fitting procedure is described as well as a non-parametric test of the existence of the log-periodic patterns for the world population. Section 4 presents a second-order generalization of the power law model, which allows for a frequency modulation in the log-periodic structure. This extended formula is used to fit the extended Dow Jones Industrial average. Section 5 summarizes what has been achieved and compares our results with previous work. In particular, we give the explicit solutions of multivariate dynamical equations for several coupled variables, such as population, technology and capital, to show that the same finite-time singularity can emerge from the interplay of these factors while each of them individually is not enough to create the singularity. Section 6 concludes by discussing a set of scenarios for mankind close to and beyond the critical time.
## 2 Singular Growth Rate
### 2.1 Tests of exponential growth
#### 2.1.1 Human population and world GDP
A faster than exponential growth is clearly observed in the human population data from year 0 up to 1970, at which the estimated annual rate of increase of the global population reached its (preliminary?) all-time peak of $`2.1\%`$. Figure 4 shows the logarithm of the estimated world population as a function of (linear) time, such that an exponential growth rate would be qualified by a linear increase. In contrast, one clearly observes a strong upward curvature characterising a “super-exponential” behaviour. A faster than exponential growth is also clearly observed in the estimated GDP (Gross Domestic Product) of the World, shown in figure 4 for the year 0 up to 2000.
#### 2.1.2 Financial indices
Over a shorter time period, a faster than exponential growth is also observed in figures 4 to 8 for a number of economic indicators such as the Dow Jones Average since the establishment of the U.S.A. in 1790 , the S&P since 1871, as well for a number of regional and global indices since 1920, including the Latin American index, the European index, the EAFE index and the World index. In all these figures, the logarithm of the index is plotted as a function of (linear) time, such that an exponential growth rate would be qualified by a linear increase. In all cases, one clearly observes in contrast a significant upward curvature characterising a “super-exponential” behaviour.
### 2.2 A first test of power law growth
#### 2.2.1 Procedure
As shown in the derivation of equation (2), it is enough that the growth rate increases with any arbitrarily small positive power of $`p(t)`$ for a finite-time singularity to develop with the characteristic power law dependence (3). Can such a behaviour explain the super-exponential behaviour documented in the figures 4-8?
The small number of data points in these time series and the presence of large fluctuations prevent the use of a direct fitting procedure with (3). Indeed, such a fit, which typically attempts to minimise the root-mean-square (r.m.s.) difference between the theoretical formula and the data, is highly degenerate: many solutions are found which differ by variations of at most a few percent of the root-mean-square (r.m.s.) of the errors. Such differences in r.m.s. are not significant, especially considering the strongly non-Gaussian nature of the fluctuations in these data sets. Maximum likelihood methods are similarly limited. To address this problem of degeneracy, we turn to a non-parametric approach, consisting in fixing $`t_c`$ and plotting the logarithm of the data as a function of $`\mathrm{log}(t_ct)`$. In such a plot, a linear behaviour qualifies the power law (3), and the slope gives the exponent $`z`$ which then can be determined visually or, better, by a fit but now with $`t_c`$ fixed. This procedure is not plagued by the previously discussed degeneracy and provides reliable and unique results.
#### 2.2.2 World population
In figures 12-12, the world population in logarithm scale is shown as a function of $`t_ct`$ also in logarithmic scale for three choices 2030, 2040 and 2050, respectively for $`t_c`$. Even though the fits with equation (3) for three cases varies in quality, they all capture the acceleration in the second half of the data on a logarithmic scale. The curvature seen in the data far from $`t_c`$ can be modeled by including a constant term in equation (3) embodying for instance the effect of an initial condition, as we discuss below. Changing $`t_c`$ from 2030 to 2050 has two competing effects observed in figures 12-12: a larger value of $`t_c`$ provides a better fit in the latter time period while deteriorating somewhat the fit to the data in the early time periods.
#### 2.2.3 World GDP
As discussed in the introduction, the human population is strongly coupled with its outputs and with the Earth’s carrying capacity, and can partly be measured by its economic production. Hence, we should expect a close relationship between the size of the human population and its GDP. Figures 12-16 show the logarithm of the estimated World GDP as a function of $`t_ct`$, both in log-log coordinates, where $`t_c`$ has been chosen to 2040, 2050 and 2060, respectively. The equation (3) is again parameterising the data quite satisfactorily.
We stress that we use the logarithm of the World GDP as well as the logarithm of the national, regional or global indices presented below as the “bare” data on which we test the power law hypothesis. This means that we plot the logarithm of the GDP or of the indices in logarithmic scale, which effectively amounts to taking the logarithm of the logarithm of the GDP as a function of $`t_ct`$, itself also in logarithmic scale in order to test for the power law (3). This is done in an attempt to minimise the effect of inflation and other systematic drifts, and in accordance with standard economic practice that only relative changes should be considered. Removing an average inflation of 4% does not change the results qualitatively but the corresponding results are not quantitatively reliable as the inflation has varied significantly over US history with quantitative impacts that are difficult to estimate.
#### 2.2.4 Financial indices
Further support for a singular power law behaviour of the economy can be found by analysing in a similar way the national, regional or global indices shown in figures 4-8. The results are shown in figures 16-32. Equation (3) is again perfectly compatible with the data and much better than any exponential model. As shown in Table 1, the fits of all six indices are found to be consistent with similar values for the exponent $`z1`$, the absolute value of the exponent increasing with $`t_c`$.
The results presented in this section on the world population, on the world GDP and on six financial indices suggest that the power law (3) is an adequate model. It is also parsimonious since the same simple mathematical expression, approximately the same critical time $`t_c`$ and same exponent are found consistently for all time series. These results confirm and extend the analysis presented forty years earlier for the world population only , which concluded at a $`t_c=2026`$. The results shown in the figures 12-12, with the sensitivity analysis provided by varying $`t_c`$ from 2030 to 2050, illustrate the large uncertainty in its determination. It is thus worthwhile to attempt quantifying further the observed power law growth and test how well $`t_c`$ is constrained.
### 2.3 Quantitative fits to a power law
In the derivation of (3), a key assumption was to neglect the limiting negative term in (1), which is warranted sufficiently close to $`t_c`$. Far from $`t_c`$, this analysis and more general considerations lead us to expect the existence of corrections to the pure power law (3). Furthermore, it may be necessary to include higher order terms as well as generalise the exponent as we will see in the next section.
The simplest extension of equation (3) is
$$p(t)=A+B(t_ct)^z.$$
(5)
In order to make a first quantitative estimate of the acceleration in the growth rate, determined by the exponent $`z`$ and the position $`t_c`$ of the singularity, we now let $`t_c`$ be a free parameter. In figure 33, the equation (5) is fitted to the world population from 0 to 1998. The parameter values of the fit are $`A0`$, $`B22120`$, $`t_c2078`$ and $`z1.9`$. The negative value of the exponent is compatible with $`A0`$. The negative exponent $`z1.9`$ obtained in the fit means that equation (5) has a singularity at $`t=t_c`$ corresponding to an infinite population. This is clearly impossible on a finite Earth. The point to be extracted from this analysis is that the world population has until very recently grown at an accelerating growth rate in good agreement with a singular behaviour. Singularities are always mathematical idealisations of natural phenomena: they are not present in reality but foreshadow an important transition or change of regime. In the present context, they must be interpreted as a kind of “critical point” signaling a fundamental and abrupt change of regime similar to what occurs in phase transitions .
As already discussed in relation with equation (1) and in the previous section, the world population growth cannot be separated from that of its evolving carrying capacity. As a first attempt to quantify this variable in an independent way, we analyse quantitatively the two largest data sets among all the financial indices and GDP: due to the large fluctuations of the financial indices compared to the number of points, only the S&P and the Dow Jones Average gave reliable results when $`t_c`$ is a free parameter. Figure 34 shows the corresponding fits with equation (5). The parameter values of the best fits are $`A14`$, $`B71`$, $`z0.27`$ and $`t_c2068`$ for the Dow Jones Average and $`A0`$, $`B1693`$, $`z1.3`$ and $`t_c2067`$ for the S&P. The fit with equation (5) exemplifies the acceleration of the growth rate, which is our main message. However, the location of the critical point is still not very reliable when based on simple power fits of very noisy data . This motivates us to extend this analyses in the following sections.
## 3 Beyond a simple power law
The results shown in the figures 12-32, with the sensitivity analysis provided by varying $`t_c`$ from 2030 or 2040 to 2050 or 2060, illustrate the large uncertainty in the determination of the critical time. The direct fit with (5) still gives a very large uncertainty. As can be seen from the figures, an important reason lies in the existence of large fluctuations around the average power law behaviour. In the next section, we will see that this variability might be genuine and not simply noise. Furthermore, adding an extra degree of freedom will certainly improve a parametrisation of the data.
### 3.1 Generalisation to power laws with complex exponents: log-periodicity
The idea is to generalise the real exponent $`z`$ to a complex exponent $`\beta +i\omega `$, such that a power law is changed into $`(t_ct)^{\beta +i\omega }`$, whose real part is $`(t_ct)^\beta \mathrm{cos}\left(\omega \mathrm{ln}(t_ct)\right)`$ . The cosine will decorate the average power law behaviour with so-called log-periodic oscillations, the name steming from the fact the oscillations are periodic in $`\mathrm{ln}(t_ct)`$ and not in $`t`$. As we shall see, these log-periodic oscillations can account for a large part of the observed variability around the power law. Thus, taking them into account provides a better parametrisation of the data and hence better constraints on the parameters of the power law $`\beta `$ and $`t_c`$.
There are fundamental reasons for introducing log-periodic corrections. Singularities often exhibit genuine log-periodic corrections that result from specific mechanisms : singularities in the Euler equations with complex exponents have been found to result from a cascade of Rayleigh-Taylor instabilities leading to log-periodic oscillatory structures around singular vortices organised according to discrete self-similar pancakes; in the process of formation of black holes, the matter field solution oscillates periodically in the logarithm of the difference between time and time of the formation of the singularity ; the phase separation kinetics of a binary mixture subjected to an uniform shear flow quenched from a disordered to a homogeneous ordered phase exhibits log-periodic oscillations due to a cyclical mechanism of stretching and break-up of domains, which allows to store and dissipate elastic energy in the system ; material failure occurs after intermittent damage acceleration and quiescent phases that are well-described by log-periodic structures decorating an overall power law singularity ; stock market crashes preceded by speculative bubbles provide an highly relevant analogy to the question of sustainability in the growth rate of the human population. More generally, from the point of view of field theory as a tool-box for constructing theories of complex systems, we should expect generically the existence of complex exponents and their associated log-periodic corrections . We suggest that the presence of log-periodic oscillations deriving from general theoretical considerations can provide a first step to account for the ubiquitous observation of cycles in population dynamics and in economics.
### 3.2 Log-periodic fit of the World population
#### 3.2.1 Results
Guided by the recent progress in the understanding of complex systems and the possibility of complex exponents discussed in the previous section, we have also fitted the world population data with the following equation
$$p(t)A_1+B_1(t_ct)^\beta +C_1(t_ct)^\beta \mathrm{cos}\left(\omega \mathrm{ln}(t_ct)+\varphi \right),$$
(6)
as shown in figure 33. We obtained two solutions, the best having $`A0`$, $`B1624`$, $`C127`$, $`z1.4`$, $`t_c2056`$, $`\omega 6.3`$ and $`\varphi 5.1`$. The second solution has $`A0.25`$, $`B1624`$, $`C127`$, $`z1.7`$, $`t_c2079`$, $`\omega 6.9`$ and $`\varphi 4.4`$. In this extension of equation (5), the cosine term embodies a discrete scale invariance decorating the overall acceleration with a geometrical scaling ratio $`\lambda =\mathrm{exp}\left(2\pi /\omega \right)`$: the local maxima of the oscillations are converging to $`t_c`$ with the geometrical ratio $`1/\lambda `$.
#### 3.2.2 Sensitivity analysis
Due to the small number of points in the population data set, the robustness of the fit with equation (6) was investigated with respect to fluctuations in the important physical parameters $`t_c`$, $`\beta `$ and $`\omega `$ . The method we used was as follows. Together with the data set (data set 1) obtained from the United Nations Population Division Department of Economic and Social Affairs (see the introduction section), which covers the period $`\left[0:1998\right]`$, seven other data sets where analysed in an identical manner. These first three data sets were generated by removing the first point (data set 2), the two first points (data set 3) and the 3 first points (data set 4). Hence, those three data sets cover the periods $`\left[1000:1998\right]`$, $`\left[1250:1998\right]`$ and $`\left[1500:1998\right]`$. A fifth data set (data set 5) was constructed by including the UN estimate that the world’s population would reach 6 billion in October 1999 to the original data set (data set 1). Three additional data sets were created by removing points in the other end from the original data set (data set 1), i.e, by removing the last point (data set 6), the two last points (data set 7) and the three last points (data set 8). Hence, those three data sets cover the periods $`\left[0:1990\right]`$, $`\left[0:1980\right]`$ and $`\left[0:1970\right]`$.
The differences between the results obtained for the first five data sets are minor, as can be seen in Table 2 showing the values corresponding to the best fits. Data set 6 and 7 are also compatible with the previous 5 whereas the fit to data set 8 exhibit a significant discrepancy. For $`t_c`$, this gives the window 2052 $`\pm `$ 10 years, which is rather well-constrained. Furthermore, the values obtained for $`\omega 6\pm 0.5`$ (again except for data set 8) are also quite compatible with previous results. The corresponding fluctuations in the fundamental parameters $`z1.35\pm 0.11`$ and $`\lambda 2.8\pm 0.3`$ are also within reasonable bounds. Note that it is difficult to obtain a better resolution in time as world population statistics in past centuries are all generated by using some sort of statistical regression model. This might explain the relatively low value of the spectral peak obtained for data set 5, see below. Furthermore, the peak clearly stands out against the background for seven out of eight spectra as we now discuss. Another encouraging observation is the notable amplitude of the log-periodic oscillations quantified by $`C`$, approximately $`510\%`$ of the pure power law acceleration quantified by $`B`$, as seen in the caption of figure 33.
#### 3.2.3 Non-parametric tests of log-periodicity
We also present a non-parametric test for the existence of the log-periodic oscillations decorating the spontaneous singularity, obtained by eliminating the leading trend using the transformation
$$p\left(t\right)\frac{p\left(t\right)A_1B_1(t_ct)^\beta }{C_1(t_ct)^\beta }.$$
(7)
This transformation should produce a pure $`\mathrm{cos}\left(\omega \mathrm{ln}(t_ct)+\varphi \right)`$ if equation (6) was a perfect description. In figure 36, we show the residual defined by (7) for data 3 and data 5 as a function of $`\mathrm{ln}(t_ct)`$ as well as their Lomb periodograms which provide a power spectrum analysis for unevenly sampled data: the approximately regular oscillations in $`\mathrm{ln}(t_ct)`$ give a significant spectral peak at a log-angular frequency $`\omega 5.86.1`$ compatible with the fit of equation (6), see Table 2.
### 3.3 Summary
To sum up the evidence obtained so far, the comparison between the semi-logarithmic plots in figures 4-8 and the log-log plots in figures 12-32 validate the power law model (3) at the expense of the exponential model: there is no doubt that the world population and major economic and financial indices on average have grown much faster than exponentially. The second message is that the rather large fluctuations decorating an average power law acceleration can be remarkably well described by a simple generalisation of the power law in terms of a complex exponent: not only do we see a good agreement between the spectral analysis and the fits with equation (6), in addition the small fluctuations in the values for $`t_c`$, $`\beta `$ and $`\omega `$ for the 7 of the 8 data sets make the analysis credible for the world population. Of course, this does not prove that equation (6) is the correct description and equation (5) is a wrong description. However, since the r.m.s. of the fits with the two equations differs by a factor of $`4`$, there is no doubt that equation (6) does a better job of parameterising the data. This is the numerical argument. The theoretical justification has already been given above. The two combined certainly makes the case stronger. For the financial indices, the use of equation (6) does not lead to a significant improvement, and this leads us to examine the relevance of the next order of the expansion of corrections to the power law.
## 4 To second order
### 4.1 Next order of the log-periodic expansion
The data set containing the Dow Jones Average consists of $`2500`$ monthly quotes for the period $`\left[1790:1999.9\right]`$. We propose that it is representative of the capitalistic growth of the U.S.A. The time span and the sampling rate of this data set makes it reasonable to use the generalisation (12) of (6) to second order which allows for a continuous shift in the angular log-frequency $`\omega `$ in what effectively corresponds to a Landau or renormalisation group expansion depending on the prefered framework.
We briefly summarize the method. Using the renormalization group (RG) formalism on a financial index $`I`$ amounts to assuming that the index at a given time $`t`$ is related to that at another time $`t^{}`$ by the transformations
$$x^{}=\varphi (x),$$
(8)
$$F(x)=g(x)+\frac{1}{\mu }F(\varphi (x)),$$
(9)
where $`x=t_ct`$. $`t_c`$ is the critical time and $`\varphi `$ is called the RG flow map. Here,
$$F(x)=I(t_c)I(t),$$
such that $`F=0`$ at the critical point and $`\mu `$ is a constant describing the scaling of the index evolution upon a rescaling of time (8). The function $`g(x)`$ represents the non-singular part of the function $`F(x)`$. We assume as usual that the function $`F(x)`$ is continuous and that $`\varphi (x)`$ is differentiable. In order to use this formalism to constrain the possible time dependence of the index, we notice that the solution in terms of a power law of the RG equation (9) together with (8) and the linear approximation $`\varphi (x)=\lambda x`$ valid close to the critical point can be rewritten as
$$\frac{dF(x)}{d\mathrm{ln}x}=\alpha F(x).$$
(10)
This states simply that a power law is nothing but a linear relationship when expressed in the variables $`\mathrm{ln}F(x)`$ and $`\mathrm{ln}x`$. A critical point is characterized by observables which have an invariant description with respect to scale transformations on $`x`$. We can exploit this and the expression (10) to propose the structure of the leading corrections to the power law with log-periodicity. Hence, we notice that (10) can be interpreted as a bifurcation equation for the variable $`F`$ as a function of a fictitious “time” ($`\mathrm{ln}x`$) as a function of the “control parameter” $`\alpha `$. When $`\alpha >0`$, $`F(x)`$ increases with $`\mathrm{ln}x`$ while it decreases for $`\alpha <0`$. The special value $`\alpha =0`$ separating the two regimes corresponds to a bifurcation. Once we have recognized the structure of the expression (10) in terms of a bifurcation, we can use the general reduction theorem telling us that the structure of the equation for $`F`$ close to the bifurcation can only take a universal non-linear form given by
$$\frac{dF(x)}{d\mathrm{ln}x}=(\alpha +i\omega )F(x)+(\eta +i\kappa )|F(x)|^2+𝒪(F^3).$$
(11)
where $`\alpha >0`$, $`\omega `$, $`\eta `$ and $`\kappa `$ are real coefficients and $`𝒪(F^3)`$ means that higher order terms are neglected. The generality of this expression stems from the fact that it is nothing but a Taylor’s expansion of a general functional form $`\frac{dF(x)}{d\mathrm{ln}x}=(F(x))`$. Such expansions are known in the physics literature as Landau expansions. We stress that this expression represents a non-trivial addition to the theory, constrained uniquely by symmetry laws. Going up to second order included, equation $`(\text{6})`$ becomes
$$\mathrm{ln}\left(I\left(t\right)\right)=A_2+B_2\frac{\left(t_ct\right)^\beta }{\sqrt{1+\left(\frac{t_ct}{\tau }\right)^{2\beta }}}[1+C_2\mathrm{cos}(\omega \mathrm{ln}(t_ct)+\frac{\mathrm{\Delta }\omega }{2\beta }\mathrm{ln}(1+\left(\frac{t_ct}{\tau }\right)^{2\beta }\varphi )\left)\right].$$
(12)
This extension has been found useful in order to account for the behaviour of stock market prices before large crashes over extended period of times up to 8 years. The present analysis thus constitutes a major generalisation as it includes over 200 years of data. Previous work have established a robust and universal signature preceding large crashes occuring in major financial stock markets, namely accelerated price increase decorated by large scale log-periodic oscillations culminating in a spontaneous singularity (critical point). The previously reported cases, which are well-described by equation $`(\text{6})`$, comprise the Oct. 1929 US crash, the Oct. 1987 world market crash, the Jan. 1994 and Oct. 1997 Hong-Kong crashes, the Aug. 1998 global market event, the April 2000 Nasdaq crash, the 1985 Forex event on the US dollar, the correction on the US dollar against the Canadian dollar and the Japanese Yen starting in Aug. 1998, as well as the bubble on the Russian market and its ensuing collapse in June 1997 . Furthermore, twenty-one significant bubbles followed by large crashes or by severe corrections in the stock markets indices of the South American and Asian countries, which exhibit log-periodic signatures decorating an average power law acceleration, have also been identified . In all these analyses, the time scales have been restricted to 1 to 8 years. In contrast, the general renormalisation group theory of such spontaneous singularities allow for an hierarchy of critical points at all scales . The results given below suggest that singularities do indeed cascade in a robust way up to the largest time scales or conversely from the largest scale to the smallest scales .
### 4.2 Second order fit of the Dow Jones Average index
#### 4.2.1 Methodology
We fit the logarithm of the extended Dow Jones index to equation (12). As mentioned, taking the logarithm provides in our opinion the simplest and most robust way to account for inflation. Furthermore, taking the logarithm embodies the notion that only relative changes are important. Another more subtle reason can be given in terms of the magnitude of the crash following the singularity: a simple model of rational expectations shows that if the loss during a crash is proportional to the maximum price, then the relevant quantity is the logarithm of the price in accordance with the standard economic notion that only relative changes should be relevant.
Fitting equation (12) to some data set is difficult even with a large data set (for noisy data with only a few hundred points or less, it becomes quite impossible), due to the degenerate r.m.s. landscape corresponding to the existence of many local minima as a function of the free parameters $`t_c`$, $`\beta `$, $`\omega `$, $`\tau `$ $`\mathrm{\Delta }\omega `$ and $`\varphi `$. This means that the r.m.s. alone is not a good measure of the quality of the fit and additional physical constraint are needed as discriminators. This has been discussed at length in . In brief, we will demand that the value of $`\beta `$ and $`\omega `$ are compatible with what has been found previously for large crashes and that the value of the transition time $`\tau `$ between the two competing frequencies is compatible with the time window $`t_ct_0`$, where $`t_0`$ is the date of the first data point and $`t_c`$ the date of the singularity in the first derivative. Unfortunately, we have no means to impose a criterion on the frequency-shift $`\mathrm{\Delta }\omega `$.
Specifically, we will demand that $`0.2<\beta <0.7`$, $`4.5<\omega <9`$ and $`\left(t_ct_0\right)/3<\tau \stackrel{<}{}\left(t_ct_0\right)`$ and the more the parameters fall in the mid-range, the higher confidence is attributed to the fit. These constraints are similar to what was used in except for the constrain on $`\tau `$ which upper limit has been made stricter here. The reason for this is simply that, whereas in the cases of the 1929 and 1987 stock market crashes on Wall Street, it was not obvious to decide the starting date of the bubble, it is now objectively determined by a historical event being the creation of the U.S.A. as an independent nation . The parameter values of the five qualifying fits is shown in Table 3. We stress that the majority of the fits were discarded due to rather large values for either $`\omega `$ or $`\tau `$ or negative values for $`\omega `$. We see that the best fit in terms of the r.m.s. also has the most reasonable parameter values for $`\beta `$, $`\omega `$ and $`\tau `$ in terms of the discussion above.
#### 4.2.2 Results
The best fit of equation (12)) to the $`210`$ years of monthly quotes is shown in figure 36 and its parameter values are given in the caption. Note that the value of the angular log-frequency $`\omega 6.5`$ compared to $`\omega 6.3`$ as well as the value for the position of the singularity $`t_c2053`$ compared $`t_c2056`$ are in close agreement with the values found for the analysis of the world population. Furthermore, the cross-over time scale $`\tau 171`$ years is perfectly compatible with the total time window of 210 years. In figure 37, the relative error between the fit and the data is shown. We see that the error fluctuates nicely around zero as it should. Furthermore, the error is decreasing from left to right clearly showing that the acceleration in the data is better and better modeled by equation (12)) as we approach the present. This behaviour is in fact to be expected from an equation such as (2) allowing for an additive noise term to describe other sources of uncertainties: using the Fokker-Planck formalism, one can show that, as the singularity at $`t_c`$ is approached, the noise term becomes negligible and the acceleration of the data should approach better and better a pure power law. This can also be seen directly from (2) with an additive noise: the divergence of $`p(t)`$ dwarves any bound noise contribution.
#### 4.2.3 Discussion
The inset of figure 36 shows the extrapolation of the fit up to the critical time $`t_c=2053`$. It suggests that the Dow Jones index will climb to impressive values in the coming decades from its present level around 11,000 at the beginning of year 2000. It is interesting that this resonates with a series of claims that the Dow Jones will climb to 36,000 , 40,000 or even 100,000 in the next two or three decades. Glassman, an investing columnist for the Washington Post, and Hassett, a former senior economist with the Federal Reserve, develop the argument that stocks have been undervalued for decades and that, for the next few years, investors can expect a dramatic one-time upward adjustment in stock prices . Elias, a financial advisor and author, believes that forces such as direct foreign investment, domestic savings, and cooperative central-banking policies will drive the vigorous market, as will the dynamics of the New Economy, which allows for the coexistence of high economic growth, low interest rates, and low inflation. In his view, the Dow Jones could reach 40,000 around 2016 . Kadlec, chief investment strategist for Seligman Advisors Inc. predicts that the Dow Jones Industrial Average will end up at 100,000 in the year 2020 . We find that equation (12) predicts that the level 36,000-40,000 will be reached in 2018-2020 A.D. and the level 100,000 in 2026 A.D, not far from these claims! Of course, the extrapolation of this growth closer to the singularity becomes unreliable due to standard limitations, such as finite size effects, and must be taken with a “hand-full of salt”.
In the academic financial literature, a time series such as the Dow Jones shown in figure 36 has been argued to exhibit an anomalously large return, averaging $`6\%`$ per year over the 1889-1978 period , which cannot be explained by any reasonable risk aversion coefficient. A solution for this puzzle is that infrequent large crashes occur or even a major still untriggered crash is looming over us; in this interpretation, the “anomalous” return becomes the normal remuneration for the risk to stay invested in the market . Our analysis suggests that the situation is even worse than this: not only the market has a large growth rate but this growth rate is accelerating such that the market is growing as a power law towards a spontaneous singularity.
## 5 Synthesis and theoretical discussion
### 5.1 Summary
The fact, that both the human world population over two thousand years, the GDP of the world and six national, regional and world financial indices over most of their lifespan agree both in i) the prediction of a spontaneous singularity, ii) the approximate location of the critical time and iii) the approximate self-similar patterns decorating the singularity is quite remarkable to say the least. This suggests that they may have a closely correlated dynamics, in fact more than the coupling between population $`p(t)`$ and carrying capacity $`K(t)`$ written in equations such as $`(\text{1})`$ would make us believe. The outstanding scientific question is whether the rate of innovations fueling the economic growth is a random process on which industrial and population selection operates or if it is driven by the pressing needs of the growing population. The main message of this study is that, whatever the answer and irrespective of one’s optimistic or pessimistic view of the world sustainability, these important pieces of data all point to the existence of an end to the present era, which will be irreversible and cannot be overcome by any novel innovation of the preceding kind, e.g., a new technology that makes the final conquest of the Oceans and the vast mineral resources there possible. This, since any new innovation is deeply embedded in the very existence of a singularity, in fact it feeds it. As a result, a future transition of mankind towards a qualitatively new level is quite possible.
The reader not familiar with critical phenomena and singularities may dismiss our approach without further ado on the basis that all demographic insights show that the population growth is now decelerating rather than accelerating. Indeed, many developed countries show a substantial reduction in fertility. However, “the tree should not hide the forest” as the proverb says, in other words this deceleration is compatible with the concept of a finite-time singularity in the presence of so-called “finite size effects” . Namely, it is well-known that nature does not have pure singularities in the mathematical sense of the term. Such critical points are always rounded off or smoothed out by the existence of friction and dissipation and by the finiteness of the system. This is a well-known feature of critical points . Finite-time singularities are similarly rounded-off by frictional effects, A clear example is provided by Euler’s disk , a rotating coin settling to rest in finite time after, in principle, an infinite number of rotations. In reality, the rotational speed accelerates until a point when friction due to air drag and solid contact with the support saturate this acceleration and stop the rotation abruptly. The upshot here is that finite-size effect and friction do not prevent the effect we document here to be present, namely the acceleration of the growth rate, up to a point where the proximity to the critical point makes finite size effects and dissipation-like effects to take over. The fact that these “imperfections” become relevant in the ultimate stage of the trajectory does not change the validity of the conclusions. The change of regime to a new phase subsists. Only its absolute abruptness is replaced by a somewhat smoother transition, albeit still rather sharp on the time scale of the total time span. In the present context, the observed very recent deceleration of the growth rate can be taken as a signature that mankind is entering in the critical region towards a transition to a new regime. Since the world population growth rate topped in 1970, this corresponds to approximately 80 years from the predicted critical point, or only $`4\%`$ of the total timespan of the investigated time series.
### 5.2 Related work
Other authors have documented a super-exponential acceleration of human activity. Kapitza has recently analysed the dynamical evolution of the human population , both aggregated and regionally and also documents a consistent overall acceleration until recent times. He introduces a saturation effect to limit the blow-up and discuss different scenarios. Using data from the Cambridge encyclopedia, he argues that epochs of characteristic evolutions or changes shrink as a geometrical series. In other words, the epoch sizes are approximately equidistant in the logarithm of the time to present. In a study of an important human activity, van Raan has found that the scientific production since the 16th century in Europe has accelerated much faster than exponentially . Using the data of DeLong , Hanson finds that the history of the world economic production since prehistoric times can only be accounted for by adding three exponentials, each one being interpreted as a new “revolution” : hunting followed by farming and then by industry. He finds that each exponential mode grew over one hundred times faster than its predecessor. He also plots the logarithm of the world product as a function of the logarithm of $`t_ct`$ with $`t_c=2050`$ and find a reasonable straight line decorated by oscillations marking the different transitions.
Macro-economic models have been developed that predict the possibility of accelerated growth . Maybe the simplest model is that of Kremer who notes that, over almost all human history, technological progress has led mainly to an increase in population rather than an increase in output per person. In his model, the economic output per person $`Y(t)/L(t)`$, where $`Y(t)`$ is the total output comprising all artifacts and $`L(t)`$ is the total population, is thus set equal to the subsistence level $`\overline{y}`$ which is assumed fixed:
$$\frac{Y(t)}{L(t)}=\overline{y}.$$
(13)
The output is supposed to depend on technology and knowledge $`A(t)`$ and labour (proportional to $`L(t)`$):
$$Y(t)=Y_0\left[A(t)L(t)\right]^{1\alpha },$$
(14)
where $`0<\alpha <1`$. The growth rate of knowledge and technology is taken proportional to population and to knowledge:
$$\frac{dA}{dt}=BL(t)A(t),$$
(15)
embodying the concept that a larger population offers more opportunities for finding exceptionally talented-people who will make important innovations and that new knowledge is obtained by leveraging existing knowledge. Eliminating $`Y(t)`$ and $`A(t)`$ between (13-15) gives the equation for the total population:
$$\frac{dL}{dt}=\frac{1\alpha }{\alpha }B[L(t)]^2.$$
(16)
This is the case $`\delta =1`$ of equation (2), showing that the population and its output develop a finite-time singularity (3) with the exponent $`z=1`$. Kremer tested this prediction by using population estimates extending back to 1 million B.C., constructed by archaeologists and anthropologists: he showed that the population growth rate is approximately linearly increasing with the population , in agreement with (16). Our result $`z1.9`$ for the human population exaggerates the singularity. On the other hand, as shown in Table 1, we find a remarkable consistent value $`z1`$ for all financial indices. Our refinements with the log-periodic formulas in order to account for the significant structures decorating the average power laws necessary lead to deviations from this “mean-field” value, which should be considered as an approximation neglecting the effect of fluctuations. This theory also predicts, in agreement with historical facts, that in the historical times when regions were separated, technological progress was faster in regions with larger population, thus explaining the differences between Eurasia-Africa, the Americas, Australia and Tasmania.
### 5.3 Multivariate finite-time singularities
Kremer’s model is only one of a general class of growth models . We briefly recall the general framework developed by Romer , which allows us to generalise the concept of finite-time singularities to multivariate dynamics and to exhibit the structure of its solution and follow in our exposition. The model involves four variables, labour $`L`$, capital $`K`$, technology $`A`$ and output $`Y`$. There are two sectors, a goods-producing sector where output is produced and an R&D sector where additions to the stock of knowledge are made. The fraction $`a_L`$ of the labour force is used in the R&D sector and the fraction $`1a_L`$ in the goods-producing sector; similarly, the fraction $`a_K`$ of the capital stock is used in R&D and the rest in goods production. Both sectors use the full stock of knowledge. The quantity of output produced at time $`t`$ is defined as
$$Y(t)=\left[(1a_K)K(t)\right]^\alpha \left[A(t)(1a_L)L(t)\right]^{1\alpha },$$
(17)
with $`0<\alpha <1`$. Expression (17) uses the so-called Cobb-Douglas functional form with power law relationships which imply constant returns to capital and labour: within a given technology, doubling the inputs doubles the amount that can be produced. Expression (17) writes that the economic output increases with invested capital, with technology and R&D and with labor.
The production of innovation is written as
$$\frac{dA}{dt}=B\left[a_KK(t)\right]^\beta \left[a_LL(t)\right]^\gamma \left[A(t)\right]^\theta ,B>0,\beta 0,\gamma 0.$$
(18)
The growth of knowledge is thus controlled by the pre-existing knowledge, by capital investment in research and by the size of the population of innovators.
As in the Solow model , the saving rate $`s`$ is exogenous and constant and depreciation is set to zero for simplicity so that
$$\frac{dK}{dt}=sY(t)=s\left[(1a_K)K(t)\right]^\alpha \left[A(t)(1a_L)L(t)\right]^{1\alpha }.$$
(19)
Let us consider (18). If $`K`$ and $`L`$ are constant, it reduces to an equation of the form (2), which exhibits a finite-time singularity only for $`\theta >1`$. In the presence of the coupling to the other growing dynamical variables $`K`$ and $`L`$, a finite-time singularity may occur even in the situation $`\theta <1`$.
As a first example, let us consider the case of a fixed population $`L(t)=`$ constant. Equations (18) and (19) can be rewritten as
$`{\displaystyle \frac{dA}{dt}}`$ $`=`$ $`bA^\theta K^\beta .`$ (20)
$`{\displaystyle \frac{dK}{dt}}`$ $`=`$ $`aA^{1\alpha }K^\alpha ,`$ (21)
We look for the condition on the exponents such that $`A(t)`$ and $`K(t)`$ exhibit a finite-time singularity. We thus look for solutions of the form
$`A(t)`$ $`=`$ $`A_0(t_ct)^\delta ,`$ (22)
$`K(t)`$ $`=`$ $`K_0(t_ct)^\kappa ,`$ (23)
with $`\delta `$ and $`\kappa `$ positive. Inserting these expressions in (21) and (20) leads to two equations for the two exponents $`\delta `$ and $`\kappa `$ obtained from the conditions that the powers of $`(t_ct)`$ are the same on the r.h.s. and l.h.s. of (21) and (20). Their solution is
$`\delta `$ $`=`$ $`{\displaystyle \frac{1+\beta \alpha }{(1\alpha )(\theta +\beta 1)}},`$ (24)
$`\kappa `$ $`=`$ $`{\displaystyle \frac{2\theta \alpha }{(1\alpha )(\theta +\beta 1)}}.`$ (25)
The condition that both $`\delta `$ and $`\kappa `$ are positive enforce that $`\theta +\beta >1`$, which is the condition replacing $`\theta >1`$ for the existence of a finite-time singularity in the monovariate case. This shows that the combined effect of past innovation and capital has the possibility of creating an explosive growth rate even when each of these factors in isolation does not. Note that inequality $`\theta +\beta >1`$ ensures that $`\delta >\kappa `$, i.e., the growth of the technological stock is faster than that of the capital.
There are many ways to reinsert the dynamical evolution of the population. Let us here consider the simplest one used by Kremer , which consists in assuming that $`L(t)`$ is proportional to $`K(t)`$ as given by (13). Then, expressions (18) and(19) give
$`{\displaystyle \frac{dA}{dt}}`$ $`=`$ $`a^{}\left[L(t)\right]^{\beta +\gamma }\left[A(t)\right]^\theta ,a^{}>0,\beta 0,\gamma 0,`$ (26)
$`{\displaystyle \frac{dL}{dt}}`$ $`=`$ $`b^{}L(t)\left[A(t)\right]^{1\alpha }.`$ (27)
Looking for solutions of the form (22) and (23) gives
$`\delta `$ $`=`$ $`{\displaystyle \frac{1}{1\alpha }},`$ (28)
$`\kappa `$ $`=`$ $`{\displaystyle \frac{2\theta \alpha }{\beta +\gamma }}.`$ (29)
It is interesting to find that the technology growth exponent $`\delta `$ is not at all controlled by $`\theta `$ nor $`\beta `$ and $`\gamma `$. This illustrates that a finite-time singularities can be created from the interplay of several growing variables resulting in a non-trivial behaviour. In the present context, it means that the interplay between different quantities, such as capital and technology, may produce an “explosion” in the population even though the individual dynamics do not. In particular, this interplay provides an explanation of our finding of the same value of the critical time $`t_c(2052\pm 10)`$ both for the population and economic indices.
## 6 Possible scenarios
We now attempt to guess what could be the possible scenarios for mankind close to and beyond the critical time $`t_c`$.
A gloomy scenario is that humanity will enter a severe recession fed by the slow death of its host (the Earth), in the spirit of the analogy proposed between the human species and cancer. This worry about human population size and growth is shared by many scientists, including the Union of Concerned Scientists (comprising 99 Nobel Prize winners) which asks nations to “stabilise population.” Representatives of national academies of science from throughout the world met in New Delhi, 24-27 October 1993, at a “Science Summit” on World Population. The participants issued a statement, signed by representatives of 58 academies on population issues related to development, notably on the determinants of fertility and concerning the effect of demographic growth on the environment and the quality of life. The statement finds that “continuing population growth poses a great risk to humanity,” and proposes a demographic goal: “In our judgment, humanity’s ability to deal successfully with its social, economic, and environmental problems will require the achievement of zero population growth within the lifetime of our children” and “Humanity is approaching a crisis point with respect to the interlocking issues of population, environment and development because the Earth is finite” . Possible scenarios involve a systematic development of terrorism and the segregation of mankind into at least two groups, a minority of wealthy communities hiding behind fortresses from the crowd of “barbarians” roaming outside, as discussed in a recent seminar at the US National Academy of Sciences. Such a scenario is also quite possible for the relation between developed and developing countries.
On a more positive note, it may be that “ecological” actions of the kind mentioned above will grow in the next decades, leading to a smooth transition towards an ecologically-integrated industry and humanity. Some signs may give indications of this path: during the 1990s, wind power has been growing at a rate of 26% a year and solar photo-voltaic power at 17% compared to the growth in coal and oil under 2%; governments have “ratified” more than 170 international environmental treaties, on everything from fishing to decertification . However, there are serious resistances , in particular because there is no consensus on the seriousness of the situation: for instance, the economist J.L. Simon writes that “almost every measure of material and environmental human welfare in the United States and in the World shows improvement rather than deterioration” . It may be that the strikingly similar explosive trend in population and GDP would not necessarily persist in the future when taking the differences between regional developments into account. Perhaps what is needed to avoid the finite-time singularity is a massive transfer of resources from developed to developing countries. The recent discussions at the G7/8 summit indicates that the developed world is becoming increasingly aware of the discrepancy.
Extrapolating further, the evolution from a growth regime to a balanced symbiosis with nature and with the Earth’s resources requires the transition to a knowledge-based society, in which knowledge, intellectual, artistic and humanistic values replace the quest for material wealth. Indeed, the main economic difference is that “knowledge” is non-rival : the use of an idea or of a piece of knowledge in one place does not prevent it from being used elsewhere; in contrast, say an item of clothing by an individual precludes its simultaneous use by someone else. Only the emphasis on non-rival goods will limit ultimately the plunder of the planet. Some so-called “primitive” societies seem to have been able to evolve into such a state .
The race for growth could however continue or even be enhanced if fundamentally new discoveries at a different level of the hierarchy witnessed until present enabled mankind to start the colonisation of other planets. The conditions for this are rather drastic, since novel modes of much faster propulsions are required as well as revolutions in our control of the adverse biological effects of space on humans. It may be that some evolved form of humans will appear who are more adapted to the hardship of space. This could lead to a new era of renewed accelerated growth after a period of consolidation, culminating in a new finite-time singularity, probably centuries in the future.
Acknowledgement: We thank P. Kendall and R. Prechter for help in providing the financial data from the Foundation For The Study Of Cycles, R. Hanson for the world GDP data and useful discussions, B. Taylor of Global Financial Data for the permission to use their data, M. Lagier, D. Zajdenweber for discussions, U. Frisch and D. Stauffer for a critical reading of the manuscript and for useful suggestions.
Note Added in Proofs: Nottale, Chaline and Grou have recently independently applied a log-periodic analysis to the main crises of different civilisation. They first noticed that historical events seem to accelerate. This was actually anticipate by Meyer who used a primitive for of log-periodic acceleration analysis . Grou has demonstrated that the economic evolution since the neolithic can be described in terms of various dominating poles which are subjected to an accelerating crisis/ no-crisis pattern. Their quantitative analysis on the median dates of the main periods of economic crisis in the history of Western civilization (as listed in are as follows (the dominating pole and the date are given in years / JC): {Neolithic: -6500}, {Egypt: -3000},{Egypt: -900}, {Grece: -100}, {Rome: +400}, {Byzance: +800}, {Arab expansion: +1100}, {Southern Europ: +1400}, {Netherland:+1650}, {Great-Britain: +1775}, {Great-Britain: +1830}, {Great-Britain: +1880}, {Great-Britain: +1935}, {United-States: +1975}. Log-periodic acceleration with scale factor $`\lambda =1.32\pm 0.018`$ occurs towards $`t_c=2080\pm 30`$. Agreement between the data and the log-periodic law is statistically highly significant ($`t_{\mathrm{student}}=145`$, Proba $`<<10^4`$). It is striking that this independent analysis based on a different data set gives a critical time which is compatible with our own estimate $`2052\pm 10`$. |
warning/0002/cond-mat0002312.html | ar5iv | text | # Linear and Nonlinear Susceptibilities of a Decoherent Two Level System
## I Introduction
A vast body of work has been devoted to understanding the transition to decoherence in models of a two-level-system (TLS) coupled to an infinite set of environmental degrees of freedom . The central quantity in this transition is the dynamical correlation function, or equivalently, the linear susceptibility, $`\chi ^{\prime \prime }(\nu )`$, of the system degrees of freedom. To our knowledge, no studies of the nonlinear susceptibility have been undertaken. Our interest in the nonlinear response is twofold: 1) experimental probes of Macroscopic Quantum Coherence (MQC) might access the nonlinear regime and 2) the nonlinear susceptibility appears to bear a different relationship to the coherence of the TLS than the linear susceptibility. A measure of coherence of the system is the quantity of spectral weight within the resonance peak of $`\chi ^{\prime \prime }(\nu )`$ associated with the TLS—this is simply proportional to the probability of absorbing a single photon. Absorption of a photon at a precise energy is a property uniquely associated with the underlying coherence of the quantum mechanical system. As the coupling to the environment is increased and the system increasingly localizes in one state, the probability of single photon absorption diminishes. In this paper we address whether two photon absorption is qualitatively similar—is it similarly linked to the underlying coherence of the system? We are thus led to the study of nonlinear susceptibilities of a TLS as the system is tuned through a decoherence transition.
N. D. Mermin investigated a model obtained from the well known spin-boson Hamiltonian by keeping only the lowest two levels of each harmonic oscillator comprising the environmental bath . The Hamiltonian therefore contains a “system” spin-1/2, represented by the Pauli matrices $`\sigma `$, coupled to a set of $`N`$ “environment” spin-1/2 degrees of freedom, $`\{𝐬_j\}`$. The resulting Hamiltonian
$$H=\frac{t}{2}\sigma _x+\frac{\lambda }{4N}\sigma _z\underset{j=1}{\overset{N}{}}(s_j^++s_j^{})+\frac{\omega }{2N}\underset{j=1}{\overset{N}{}}s_j^z$$
(1)
(henceforth referred to as the Mermin Model) can be solved to demonstrate the close correspondence between decoherence and a second order phase transition. The transition in the Mermin model also bears some similarity to the decoherence transition in the Ohmic case of the spin-boson Hamiltonian . Owing to its finite Hilbert space, exact dynamical susceptibilities (linear and nonlinear) can be calculated and followed through the transition. The finite dimensionality may also be realistic in some physical settings and exhibits the expected corrections to an “infinite environment” thermodynamic phase transition. In this paper, we will present such calculations and a careful study of the Mermin Model at the linear level, demonstrating an analogy to the superconducting transition in the “dirty” limit. We will then go on to the calculation and interpretation of the third order, nonlinear susceptibility corresponding to two photon absorption.
To briefly summarize our findings: The transition to decoherence is induced by strong coupling to environmental degrees of freedom (large $`\lambda `$) or a small system energy scale (small $`t`$). Below a critical coupling, $`\lambda _c`$, all spectral weight of the dynamical susceptibility
$$\chi ^{\prime \prime }(\nu )\text{Im}\frac{i}{4}_{\mathrm{}}^{\mathrm{}}𝑑t0|[\sigma (t),\sigma (0)]|0\theta (t)e^{i\nu t}$$
(2)
lies in the principal resonance of the TLS at an energy $`t`$. When the coupling is increased above the critical coupling, a new exponentially small energy scale $`O(te^{N/2})`$ emerges, associated with a broken symmetry, $`\sigma 0`$, in the thermodynamic limit $`N\mathrm{}`$. This feature corresponds to tunneling modified by a Franck-Condon type overlap factor. As the coupling is increased, spectral weight is continuously shifted to the “near-zero” frequency channel. The weight of the delta-function, $`\delta (0^+)`$, is simply the order parameter, $`|0|\sigma |0|^2`$ (or, more exactly, $`|1|\sigma |0|^2`$, for large but finite $`N`$). However, the dynamical susceptibility obeys a sum rule implying that an incompletely formed broken symmetry state leaves some spectral weight at the position of the principal resonance, set by $`t`$.
It is in this respect that the Mermin model differs from the ohmic case of the spin-boson Hamiltonian (SBH). In the SBH, the Toulouse limit (at $`\alpha =1/2`$), corresponding to complete inelastic broadening of the resonance, is a precursor to localization of the system spin (at $`\alpha =1`$) . In the Mermin model, the quantum resonance of the system spin—although damped—remains intact through the decoherence transition. This behavior is analogous to the Glover-Tinkham-Ferrell sum rule for the dynamical conductivity in superconductors. Spectral weight falling in the gap region, proportional to the superfluid density (or, equivalently, the order parameter) is redistributed to a delta-function at zero frequency. In the case of the transition to decoherence, finite but supercritical coupling, $`\lambda >\lambda _c`$, is analogous to the transition to superconductivity in the dirty limit where the superconductor can absorb energy above the gap.
In contrast, the nonlinear response has a more complicated behavior. Two photon absorption (TPA) requires the presence of ancillary levels as intermediate states and is therefore enhanced when there is some degree of “inelastic” broadening of the primary resonance. At the two extremes, TPA is zero when $`\kappa >\kappa _c`$ and the environment is decoupled from the system—but it is also zero when $`\kappa 0`$ and the system is completely decoherent.
We now turn to a review of the analytical results obtained from the Mermin model.
## II The Mermin Model
In the absence of coupling to the environment, the two level system in (1) possesses a ground state in a coherent superposition of spin up and spin down (denoted $`|+|`$) and will display a sharp resonance in $`\chi ^{\prime \prime }(\nu )`$ at $`\nu =t`$. When environmental coupling is included, there are two effects commonly associated with decoherence. First, the coherence features should shift to lower frequency (Franck-Condon overlap) and secondly, the peak should broaden from “inelastic” energy exchange with the environment. These two effects are nicely demonstrated in the Mermin Model. In (1), the environment spins may be summed to one big $`O(N)`$ spin, $`𝐒`$, and the Hamiltonian becomes
$$H=\frac{t}{2}\sigma _x+\frac{\lambda }{2N}\sigma _zS_x+\frac{\omega }{2N}S_z$$
(3)
As Mermin points out, the advantage of the Hamiltonian (3) is that in the limit $`N\mathrm{}`$, the environment spins may be replaced in the Hamiltonian by the $`x`$ and $`z`$ components of a classical spin angular momentum: $`\frac{\lambda }{N}\sigma _zS_x\frac{1}{2}\lambda \sigma _z\mathrm{sin}\theta `$ and $`\frac{\omega }{N}S_z\frac{1}{2}\omega \mathrm{cos}\theta `$. The resulting Hamiltonian
$$H=\frac{t}{2}\sigma _x+\frac{\lambda }{4}\sigma _z\mathrm{sin}\theta \frac{\omega }{4}\mathrm{cos}\theta $$
(4)
may be diagonalized and its ground state eigenvalue, $`E_0(\theta )`$, given by
$$E_0(\theta )=\frac{1}{2}\sqrt{t^2+\frac{\lambda ^2}{4}\mathrm{sin}^2\theta }\frac{\omega }{4}\mathrm{cos}\theta $$
(5)
minimized with respect to $`\theta `$.
The critical behavior of the model may now be seen by examining the ground state energy $`E_0`$, as function of $`t`$. $`E_0`$ bifurcates at a finite value of $`t`$ going from a singlet, non-degenerate root for $`t>\lambda ^2/2\omega `$ to a doubly degenerate set for $`t<\lambda ^2/2\omega `$.
$`\theta _0=0`$ $`t>\lambda ^2/2\omega `$ (6)
$`\mathrm{sin}\theta _0=\pm \sqrt{{\displaystyle \frac{14t^2\omega ^2/\lambda ^4}{1+\omega ^2/\lambda ^2}}}`$ $`t<\lambda ^2/2\omega `$ (7)
In the former case, the environment is decoupled from the system and always points along the $`z`$-axis to minimize its “Zeeman” energy corresponding to the last term of (3). In the latter case, the environment and the system spins are frozen in two distinct orientations with degenerate energies. The ground state wave function is still in a superposition $`\alpha |+\beta |`$ but it is no longer an evenly weighted one ($`\alpha =\beta `$); rather, the two roots correspond to the system predominantly in $`|+`$ or $`|`$.
However, for finite (even very large) $`N`$, the system must be in an evenly weighted superposition—independent of the value of $`t`$. For $`t0`$, it is possible to estimate the tunnel splitting between the ground and first excited states. Following , we consider the adiabatic environment (or small $`t`$) limit in which the environment instantaneously adjusts to the localization of the system in a given state and thus “points” at the angle $`\theta _0`$, as shown in Fig. 1. For the system to jump to the other state, the environment must be caught in a fluctuation which points along the other direction, $`\theta _0`$. The amplitude for $`N`$ spin-1/2’s to be found at an angle $`2\theta _0`$ away is $`\mathrm{cos}^N\theta _0`$. The new tunnel splitting is thus reduced by a Franck-Condon type overlap factor:
$$t^{}=t\theta _0|R_y(2\theta _0)|\theta _0=t\mathrm{cos}^N\theta _0=t\frac{1}{(1+\lambda ^2/\omega ^2)^{N/2}}$$
(8)
where $`R_y(2\theta _0)`$ is the rotation operator. Note that this is still a purely conservative process; the resonance is perfectly sharp although at an exponentially reduced energy scale. We will refer to this resonance as the Franck-Condon resonance.
Although, for finite $`N`$ and $`t`$, the symmetry is unbroken and the ground state is an evenly weighted superposition, there must be some symptom of the crisis at $`t\lambda ^2/2\omega `$ when $`N`$ becomes large. Presumably, this symptom should be the inelastic broadening of the resonance to an over damped state. Following along the lines of the Fermi Golden Rule (FGR) calculation in , we treat the last two terms of (3) as the unperturbed Hamiltonian, $`H_0`$. $`H_0`$ is a Zeeman-like Hamiltonian with a magnetic field that depends upon the $`z`$-component of the system spin. Therefore it should be diagonalized with a $`\sigma _z`$ dependent rotation about the $`y`$-axis:
$$RR_y(\sigma _z\theta _0)=e^{i\sigma _z\theta _0S_y}$$
(9)
The transformed Hamiltonian is then
$$H_0^{}=R^{}H_0R=\frac{\omega }{2N}S_z\mathrm{cos}\theta _0(1+\frac{\lambda ^2}{\omega ^2})$$
(10)
where $`\mathrm{tan}\theta _0=\lambda /\omega `$ has been chosen to eliminate the $`S_x`$ term in $`H_0^{}`$. Now turning to the first term of (3), the rotation yields:
$$H_1^{}\frac{t}{2}R^{}\sigma _xR=\frac{t}{2}(\sigma ^+e^{i\sigma _z\theta _0S_y}+\sigma ^{}e^{+i\sigma _z\theta _0S_y})$$
(11)
Now the system is prepared in a state $`|+`$ with the environment in its ground state $`S_z=N/2`$ (which will be denoted $`|N/2`$). The transition rate is now calculated by applying Fermi’s Golden Rule, treating $`H_1^{}`$ as the perturbation.
$$\mathrm{\Gamma }=\frac{t^2}{4}\underset{m=N/2}{\overset{N/2}{}}||m|\sigma ^{}e^{+i\sigma _z\theta _0S_y}|N/2|+|^2\delta (E_fE_i)$$
(12)
As depicted in Figs. 2(a-c), $`H_1^{}`$ induces a rotation of the environment by $`\theta _0`$ about the $`y`$-axis as the system flips from $`|+`$ to $`|`$. Since the environment is now rotated away from the environment ground state of the final configuration ($`|N/2|`$) as well (Fig. 2c), the dominant contribution to $`\mathrm{\Gamma }`$ comes from an excited (rotated) final state. Looking at the geometry of Fig. 2c, the final state, $`|m=N/2\mathrm{cos}\theta _0|`$, contributes most strongly. Using the Hamiltonian $`H_0^{}`$, the energy difference is computed:
$$\mathrm{\Delta }E=E_fE_i\frac{\omega }{4}(1\mathrm{cos}\theta _0)\mathrm{cos}\theta _0(1+\lambda ^2/\omega ^2)\frac{\lambda ^2}{8\omega }$$
(13)
This calculation suggests an interpretation of the earlier result (6) for the critical coupling obtained in the $`N\mathrm{}`$ limit. When the width of the resonance $`\mathrm{\Delta }E`$ becomes comparable to $`t`$, the resonance is over damped and quantum coherence is destroyed.
Although both the adiabatic computation and the FGR one are small $`t`$ perturbation theories and involve the same environment overlap amplitude, their interpretations are quite different. In the former, the system in the $`|+`$ state leaves an “imprint” upon the environment which, because it is nearly orthogonal to the imprint left by system state $`|`$, reduces the tunneling amplitude by an exponentially small factor. The resonance must therefore shift from $`O(t)`$ to a smaller energy scale $`O(te^{N/2})`$. In the latter case, the system forces the environment to dissipate energy and the resonance is broadened by $`O(\lambda ^2/\omega )`$ but remains nominally at $`O(t)`$.
One might speculate upon two possible scenarios in which these two perturbative calculations are blended, as suggested in Figs. 3 and 4. In Fig. 3, a single resonance broadens and shifts to lower energy, decoherence occurring when the resonance becomes broadened beyond its energy scale. In this scenario, there is no explicit Franck-Condon energy scale, just the appearance of increasing spectral weight close to zero as the resonance becomes over-damped. In the other scenario, demonstrated in Fig. 4, spectral weight is only transferred to an exponentially small energy below some critical coupling; the resonance at $`t`$, although broadened, is not critically damped and remains distinct from the Franck-Condon resonance. To put it another way, the systems retains a gap of $`O(t)`$. Furthermore, even at finite $`N`$, there is essentially no “Franck-Condon effect” for $`\kappa >\kappa _c`$ (see Fig. 5). It is the second scenario that is manifested in the Mermin Model.
## III Calculations of the Linear Susceptibility: $`\chi ^{\prime \prime }(\nu )`$
We consider the Hamiltonian (1) with an additional coupling to an external field $`h(t)=h\mathrm{cos}\nu t`$:
$$H_{\text{ext}}=h\frac{\sigma _z}{2}\mathrm{cos}\nu t$$
(14)
For calculating nonlinear susceptibilities it will be necessary to use explicitly real driving fields, so we adopt explicitly real notation for the expectation value of the time dependent system spin, at this point. The time dependence of $`\sigma _z`$, at linear response, is
$$\frac{1}{2}\sigma _z(t)=h\chi ^{}(\nu )\mathrm{cos}\nu th\chi ^{\prime \prime }(\nu )\mathrm{sin}\nu t$$
(15)
At zero temperature, the response function $`\chi (\nu )`$ is given by the integral on the right hand side of (2). Expanded out,
$$\chi (\nu )=\frac{1}{\mathrm{}}\underset{n}{}\frac{1}{4}|\sigma _{0n}|^2(\frac{1}{\nu +\nu _{n0}i\delta }\frac{1}{\nu \nu _{n0}i\delta })$$
(16)
where $`\nu _{ij}(E_iE_j)/\mathrm{}`$ and $`\delta =0^+`$. The matrix elements of the system spin $`\sigma _{ij}i|\sigma _z|j`$ will be referred to as the order parameter matrix. (The state $`|i`$ now refers to the combined system and environment state.) Using (15), the rate of energy absorption is then
$$\mathrm{\Omega }=\frac{1}{2}\nu \chi ^{\prime \prime }(\nu )h^2$$
(17)
Equations analogous to (15), (16) and (17) for the third order nonlinear susceptibility are derived in the Appendix.
The Hamiltonian (1) was diagonalized and the complete order parameter matrix, $`\sigma _{ij}`$, necessary for the linear and nonlinear calculations, was computed. It is natural to work with one dimensionless coupling constant, $`\kappa 2\omega t/\lambda ^2`$. The critical coupling implied in the $`N\mathrm{}`$ calculation is then $`\kappa =\kappa _c=1`$.
Fig. 4 shows the behavior of $`\chi ^{\prime \prime }(\nu )`$ as the system energy scale $`t`$ is reduced, corresponding to a range $`\kappa =0.6\kappa _c1.2\kappa _c`$. These computations were performed for $`80`$ environment spins. Computations with up to $`300`$ environment spins suggest that there is little change beyond $`80`$ spins. An artificial broadening ($`\delta =0.01`$) was added to make the features more visible. As seen in this sequence, $`\kappa =\kappa _c`$ is marked by the appearance of the Franck-Condon resonance at an exponentially small energy scale ($`\nu _{10}10^3`$). The spectral weight remaining at $`O(t)`$ when $`\kappa >\kappa _c`$ clearly exhibits inelastic broadening, although the resonance essentially disappears before becoming critically damped.
The counterpart, at finite $`N`$, to $`\sigma 0`$ at infinite $`N`$ is the spectral weight at the Franck-Condon resonance, $`\sigma _{10}0`$. Fig. 5 demonstrates that $`\sigma _{10}`$ behaves as expected in a second order phase transition; the slope at $`\kappa =\kappa _c^+`$ grows with increasing $`N`$. The matrix, $`\sigma _{ij}`$, and consequently, $`\chi ^{\prime \prime }(\nu )`$, obey a sum rule:
$$\underset{i}{}|\sigma _{ij}|^2=1$$
(18)
An incompletely developed broken symmetry ($`|\sigma _{10}|^2<1`$) means that spectral weight remains at the principal resonance at $`O(t)`$. This behavior is somewhat analogous to the superconducting transition and the Glover-Tinkham-Ferrell sum rule for the real part of the dynamical conductivity, $`\sigma (\omega )`$:
$$_0^{\mathrm{}}\sigma _\text{n}(\omega )𝑑\omega =\frac{\pi n_\text{s}e^2}{2m}+_\mathrm{\Delta }^{\mathrm{}}\sigma _\text{s}(\omega )𝑑\omega $$
(19)
(Subscripts “n” and “s” refer to “normal” and “superconducting”.) The sum rule states that spectral weight missing from dynamical conductivity above the gap must reappear as superfluid (DC) Drude weight. This behavior may be compared with other models and computations of environmental spin decoherence . These models differ from the Mermin Model in two respects: 1) typically the environmental coupling constant is comparable to the system energy scale instead of weak $`O(\lambda /N)`$ coupling, as in the present case; 2) the coupling between the system and environmental spins is isotropic, of the form, $`\stackrel{}{\sigma }\stackrel{}{s}_j`$. These differences lead to qualitatively different decoherence effects.
## IV Nonlinear Dynamical Susceptibility
In studying the nonlinear susceptibilities, we are interested in power absorption, as this is the relevant experimental probe for MQC. Consider the response of just a single two-level-system (TLS) coupled to ancillary levels: Driving the system with a signal at frequency $`\nu `$ results in third order nonlinear responses at $`\nu `$ and $`3\nu `$. The signal at $`3\nu `$ is not absorptive and, in any case, would not be detected in any experimental setting where pumping and detection take place at (nearly) the same frequency. If the nominal frequency scale of the TLS is $`\nu _0`$, there is an imaginary but emissive response at $`\nu _0`$. (The non vanishing, odd order, nonlinear corrections to power absorption at $`\nu _0`$ must alternate in sign to produce finite absorption at $`\nu _0`$, consistent with the exact Rabi solution.) On the other hand, third order two-photon absorption (TPA) process at $`\nu _0/2`$ would resemble ordinary linear absorption of a $`\nu _0/2`$ TLS in that it would be detected as an absorptive phase shift when the pumping frequency reaches $`\nu _0/2`$. The reader is referred to ref. for further details about nonlinear processes.
Fig. 6 shows the two photon absorption spectrum along with the linear susceptibility for several values of $`\kappa `$. Our basic result is the following: as the critical value of $`\kappa `$ is reached and the two-level-system undergoes a transition to decoherence, the two photon absorption power, $`\mathrm{\Omega }_{2\gamma }`$, begins to increase. It reaches a maximum at approximately $`\kappa =\kappa _c/2`$ and then decreases to zero as $`\kappa `$ goes to zero. Fig. 7 shows the behavior of the TPA signal as a function of the control parameter $`\kappa `$.
The two photon absorption spectrum was calculated from an exact diagonalization of the Mermin model using the following expression for TPA (which is derived in the Appendix):
$`\mathrm{\Omega }_{2\gamma }=`$ $`{\displaystyle \frac{1}{2}}h^4\nu \chi ^{[3]}(\nu )`$ (20)
$`=`$ $`{\displaystyle \frac{1}{8}}({\displaystyle \frac{h}{2}})^4\nu {\displaystyle \frac{1}{\mathrm{}^3}}{\displaystyle \underset{nml}{}}\sigma _{0n}\sigma _{nm}\sigma _{ml}\sigma _{l0}`$ (22)
$`\times P_{0m}^{\prime \prime }(2\nu )P_{l0}^{}(\nu )P_{n0}^{}(\nu )`$
In the expression above, the energy denominators have been abbreviated,
$$P_{ij}(\nu )\frac{1}{\nu \nu _{ij}i\delta }$$
(23)
To calculate (20) it is necessary to find all possible sets of four couplings that connect the ground state back to the ground state. These are referred to as “four link chains”; a few such chains are depicted in Fig. 8.
To understand the non-monotonic behavior of the TPA signal, we look at the expression (20). Isolating the absorptive pole, $`\chi ^{[3]}(\nu )`$ can be factored in the following way:
$$\chi ^{[3]}(\nu )=\frac{1}{4}(\frac{1}{2})^4\frac{1}{\mathrm{}^3}\underset{m}{}P_{0m}^{\prime \prime }(2\nu )|\alpha _m(\nu )|^2$$
(24)
where $`\alpha _m(\nu )`$ is a sum over intermediate states,
$$\alpha _m(\nu )\underset{n}{}\sigma _{0n}\sigma _{nm}P_{n0}^{}(\nu )=\underset{n}{}\frac{\sigma _{0n}\sigma _{nm}}{\nu \nu _{n0}}$$
(25)
$`\alpha _m(\nu )`$ is known in nonlinear optics literature as the degenerate hyperpolarizability. We focus our attention upon the integrated spectral weight, shown below for both linear and nonlinear susceptibilities $`(\mathrm{}=1)`$:
$$I_{\text{lin}}=_{0^+}^{\mathrm{}}\chi ^{\prime \prime }(\nu )𝑑\nu I_{\text{nonlin}}=h^2_{0^+}^{\mathrm{}}\chi ^{[3]}(\nu )𝑑\nu $$
(26)
The limits of integration are meant to restrict the integration to the principal resonance and exclude the contributions from the near degenerate ground states. (These are the states $`m=0,1`$ excluded in the sum below.) The integrated weight is in some sense a “figure of merit” for detection of an MQC resonance. A meaningful comparison between linear and nonlinear absorption requires the multiplicative factor $`h^2`$ in the nonlinear expression. In both cases, the integration collapses the absorptive pole leaving, in the linear case,
$$I_{\text{lin}}=\pi (\frac{1}{2})^2\underset{n1}{}|\sigma _{0n}|^2=\pi (\frac{1}{2})^2(1\sigma _{01}^2)$$
(27)
and in the nonlinear case,
$$I_{\text{nonlin}}=\pi h^2(\frac{1}{2})^4\frac{1}{8}\underset{m0,1}{}|\alpha _m(\nu _{m0})|^2$$
(28)
Since $`I_{\text{lin}}`$ given in equation (27) depends only upon the order parameter $`\sigma _{01}`$, it has a universal behavior in terms of $`\kappa `$ through the decoherence transition (see equations (39), (40)). It turns out to be possible to express $`I_{\text{nonlin}}`$ in a similar fashion within a simple approximation and understand the non-monotonic behavior seen in fig. 7. It will be shown that $`I_{\text{nonlin}}`$ depends upon the order parameter through:
$$I_{\text{nonlin}}=\frac{2\pi h^2}{\nu _0^2}(\frac{1}{2})^4\sigma _{01}^2(1\sigma _{01}^2)+O(\frac{h^2ϵ}{\nu _0^3})$$
(29)
and thus reaches a maximum at approximately $`\kappa =\kappa _c/2`$. (In the exact solution, $`|0|\sigma _z|0|^21\kappa ^2`$ for $`\kappa _c=1`$; see equation (39).)
To understand this result, consider the usual conditions through which a broken symmetry is accessed mathematically: Apply a small static bias field $`h_{\mathrm{dc}}`$ coupled to the system through $`H_{\mathrm{dc}}h_{\mathrm{dc}}\sigma _z`$ and take the limits in the order,
$$\underset{h_{\mathrm{dc}}0}{lim}\underset{N\mathrm{}}{lim}0|\sigma _z|0$$
(30)
In the decoherent phase ($`\kappa <\kappa _c`$), a field $`h_{\mathrm{dc}}O(te^{N/2})`$ will localize the system yielding $`0|\sigma |00`$, but will have essentially no effect in the coherent phase ($`\kappa >\kappa _c`$). The relevant energy levels now look like the right half of Fig. 8 (the quasi-degenerate ground state levels are no longer coupled through $`\sigma _z`$).
The key observation is that the net transition amplitude (proportional to $`\alpha _m(\nu )`$) to go from $`|0`$ to a state $`|m`$ through degenerate intermediate states $`|n`$, $`|l`$, etc. is zero unless there is a broken symmetry leaving $`0|\sigma _z|00`$. If the intermediate levels are degenerate, $`\alpha _m(\nu )`$ may be written
$$\alpha _m(\nu )=\underset{n0,m}{}\frac{\sigma _{0n}\sigma _{nm}}{\nu \nu _{n0}}=\frac{1}{\nu \nu _{n0}}\underset{n0,m}{}\sigma _{0n}\sigma _{nm}=$$
$$=\frac{1}{\nu \nu _{n0}}(0|\sigma _z|00|\sigma _z|m+m|\sigma _z|m0|\sigma _z|m)$$
(31)
where the last line is obtained by inserting a complete set of states into $`0|\sigma _z^2|m=0`$. Since the expectation value of $`\sigma _z`$ is $`O(1)`$ only in the ground state (in the broken symmetry phase), we arrive at
$$\alpha _m(\nu )=\frac{0|\sigma _z|00|\sigma _z|m}{\nu \nu _{n0}}=\frac{\sigma _{00}\sigma _{0m}}{\nu \nu _{n0}}$$
(32)
Unless $`\sigma _{00}`$ is non zero, the transition amplitudes for $`|0|m`$ involving an intermediate state interfere destructively. Finally, the integrated weight is obtained from (28),
$`I_{\text{nonlin}}=`$ $`\mathrm{const}\times {\displaystyle \underset{m0}{}}|\alpha _m(\nu _{m0})|^2`$ (33)
$`=`$ $`{\displaystyle \frac{\mathrm{const}}{(\nu \nu _{n0})^2}}\times {\displaystyle \underset{m0}{}}\sigma _{00}^2\sigma _{0m}\sigma _{m0}`$ (34)
$`=`$ $`{\displaystyle \frac{\mathrm{const}}{(\nu \nu _{n0})^2}}\times \sigma _{00}^2(1\sigma _{00}^2)`$ (35)
where the sum rule (18) (with $`j=0`$) was used in the last line. This is the argument used to establish the result (29).
We now fill in the details of this argument appropriate for the numerical computation performed in this work by 1) setting $`h_{\mathrm{dc}}=0`$ and 2) including the finite width of the resonance. $`\alpha _m(\nu )`$ is to be evaluated at the frequency for TPA, $`\nu =\nu _{m0}/2`$—this frequency is about half of the typical frequencies of the intermediate levels $`\nu _{n0}`$. $`\alpha `$ is rewritten in such a way as to separate the quasi-degenerate ground states from the excited states at the energy of the principal resonance at $`O(t)`$ (see Fig. 9):
$$\alpha _m(\nu _{m0}/2)=\frac{\sigma _{01}\sigma _{1m}}{\nu _{m0}/2\nu _{10}}+\underset{n0,1}{}\frac{\sigma _{0n}\sigma _{nm}}{\nu _{m0}/2\nu _{n0}}$$
(36)
$`\nu _{10}=t^{}O(te^{N/2})`$ is the energy scale of the Franck-Condon resonance and will be set to zero. The difference between the resonant frequency, $`\nu _{m0}`$, and the intermediate state level, $`\nu _{n0}`$, which appears in the second term of (36) is bounded by the width of the primary resonance and is therefore typically much smaller than the primary resonance itself. Denoting a maximum value for that difference by $`ϵ`$ (that is, $`\nu _{n0}=\nu _{m0}+ϵ`$), the second term of (36) becomes:
$$\underset{n0,1}{}\frac{\sigma _{0n}\sigma _{nm}}{\nu _{m0}/2}+O(\frac{ϵ}{\nu _0^2})$$
(37)
where $`\nu _0`$ is the nominal energy scale of the resonance. Following the previous calculation for $`h_{\mathrm{dc}}0`$, the hyperpolarizibility $`\alpha `$ is now
$$\alpha _m(\nu _0)=\frac{4\sigma _{01}\sigma _{1m}}{\nu _0}+O(\frac{ϵ}{\nu _0^2})$$
(38)
Substituting $`\alpha `$ back into equation (28) and using the sum rule (18) (with $`j=1`$) establishes the result (29).
To express $`I_{\text{lin}}`$ and $`I_{\text{nonlin}}`$ in terms of the control parameter $`\kappa `$, we appeal to the exact solution in the $`N\mathrm{}`$ limit. Specifically, the matrix element $`\sigma _{01}`$ is related to the expectation value of $`\sigma _z`$ in the either of the degenerate ground states of (4) (for $`\kappa <\kappa _c`$) by:
$$\underset{N\mathrm{}}{lim}|\sigma _{01}|=|\sigma _z_0|=\sqrt{\frac{1\kappa ^2}{1+\frac{\lambda ^2}{\omega ^2}\kappa ^2}}=\sqrt{\frac{1\kappa ^2}{1+\frac{2t}{\omega }\kappa ^2}}$$
(39)
Therefore $`I_{\text{lin}}`$ and $`I_{\text{nonlin}}`$ may be compactly written in terms of $`\kappa `$:
$$I_{\text{lin}}=\pi (\frac{1}{2})^2(1\frac{1\kappa ^2}{1+\frac{\lambda ^2}{\omega ^2}\kappa ^2})$$
(40)
$$I_{\text{nonlin}}=\frac{2\pi h^2}{\nu _0^2}(\frac{1}{2})^4\frac{1\kappa ^2}{1+\frac{\lambda ^2}{\omega ^2}\kappa ^2}(1\frac{1\kappa ^2}{1+\frac{\lambda ^2}{\omega ^2}\kappa ^2})$$
(41)
Returning to fig. 7, the non-monotonic behavior of $`I_{\text{nonlin}}`$ as a function of $`\kappa `$ follows from its dependence upon the combination: $`\sigma _z^2(1\sigma _z^2)`$. A plot of $`I_{\text{nonlin}}`$ following from the analytic expression (41) is shown for comparison. In the numerical computations, the maximum intensity appears at approximately $`\kappa =1/2`$ in agreement with theory.
## V Conclusion
In this work we have attempted to make more explicit the relationship between the decoherence transition in the Mermin Model and a typical thermodynamic phase transition. At the onset of decoherence, spectral weight is transferred from above the “gap,” set by the system energy scale $`t`$, to a delta function at $`\nu =0`$. This behavior is in contrast to the SBH where symmetry breaking is prefaced, first by the complete disappearance of the inelastic peak in $`\chi ^{\prime \prime }(\nu )`$ at finite frequency (when $`\alpha =1/2`$), and then finally by the softening of renormalized system energy scale to zero (when $`\alpha =1`$).
Also, we have studied nonlinear absorption in a decoherent TLS. Single photon absorption is a measure of coherence of the system because it samples the spectral weight remaining at the principal resonance—or, the complement of the order parameter according to the sum rule. On the other hand, two-photon-absorption (TPA) samples the product of the order parameter and its complement and thus vanishes in the fully coherent or fully decoherent phase. TPA reaches maximum intensity somewhere between these two extremes.
Observing TPA may be possible in systems demonstrating quantum coherence such as the recently studied Fe<sub>8</sub> molecular nanomagnet. . In principal, measurement of the ratio of $`I_{\text{nonlin}}`$ to $`I_{\text{lin}}`$ would provide a direct measure of the order parameter $`0|\sigma _z|0`$, the localized “weight” of the system. If we consider a sample in which the system spins are individually partially localized but uncorrelated with each other (like a spin glass), TPA may be a useful probe for this type of order.
###### Acknowledgements.
The author gratefully acknowledges the support of the Cottrell Foundation through Research Corporation Grant number CC3834. Work performed at BNL supported by the U.S. DOE under contract no. DE-AC02 98CH10886. The author also wishes to acknowledge many useful discussions with Dr. Jeffrey Clayhold, Dr. Barry Friedman, Dr. V. N. Muthukumar and the help of Joshua Eisner in performing some preliminary computations.
## A Derivation of the Nonlinear Dynamical Susceptibility
We begin by expressing the full time dependence of the expectation value of the desired operator in the presence of the driving field.
$$\sigma _z(t)=0|U^{}(t,\mathrm{})\sigma _z(t)U(t,\mathrm{})|0$$
(A1)
where
$`U(t,\mathrm{})=`$ $`\text{T}e^{i_{\mathrm{}}^tH(t^{})𝑑t^{}}`$ (A2)
$`=`$ $`1\underset{I_1}{\underset{}{i{\displaystyle _{\mathrm{}}^t}𝑑t_1H(t_1)}}`$ (A6)
$`\underset{I_2}{\underset{}{{\displaystyle _{\mathrm{}}^t}𝑑t_1{\displaystyle _{\mathrm{}}^{t_1}}𝑑t_2H(t_1)H(t_2)}}`$
$`\underset{I_3}{\underset{}{+i{\displaystyle _{\mathrm{}}^t}𝑑t_1{\displaystyle _{\mathrm{}}^{t_1}}𝑑t_2{\displaystyle _{\mathrm{}}^{t_2}}𝑑t_3H(t_1)H(t_2)H(t_3)}}`$
$`+\mathrm{}`$
The time dependent operators appearing in these expressions are interaction picture operators, treating the full Hamiltonian (1) as the bare Hamiltonian and the external coupling Hamiltonian , $`H(t)=\frac{h}{2}\sigma _z(t)\mathrm{cos}\nu t`$, as the perturbation. The $`S`$-matrix is expanded to the appropriate order. Grouping the third order terms in (A1) together, we arrive at:
$``$ $`0|\sigma (t)I_3|0+0|I_3^{}\sigma (t)|0=`$ (A10)
$`{\displaystyle \frac{ih^3}{\mathrm{}^3}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t_1{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t_2{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t_30|\sigma (t)\sigma (t_1)\sigma (t_2)\sigma (t_3)|0`$
$`\times \mathrm{cos}\nu t_1\mathrm{cos}\nu t_2\mathrm{cos}\nu t_3\theta (tt_1)\theta (t_1t_2)\theta (t_2t_3)`$
$`+\text{c.c.}`$
and
$``$ $`0|I_1^{}\sigma (t)I_2|0+0|I_2^{}\sigma (t)I_1|0=`$ (A14)
$`{\displaystyle \frac{+ih^3}{\mathrm{}^3}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t_1{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t_2{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t_30|\sigma (t_1)\sigma (t)\sigma (t_2)\sigma (t_3)|0`$
$`\times \mathrm{cos}\nu t_1\mathrm{cos}\nu t_2\mathrm{cos}\nu t_3\theta (tt_1)\theta (t_1t_2)\theta (t_2t_3)`$
$`+\text{c.c.}`$
We use the integral representation of the step function,
$$\theta (t)=\frac{1}{2\pi i}_{\mathrm{}}^{\mathrm{}}𝑑\omega \frac{e^{i\omega t}}{\omega i\delta }$$
($`\delta =0^+`$) and collapse the time integrations. From the various products of delta functions, terms are grouped according to whether they oscillate at $`\pm 3\nu `$ (third harmonic generation) or $`\pm \nu `$. From the latter terms, absorptive contributions come from the part of the signal that is $`\pi /2`$ out of phase with the driving term, $`\frac{h}{2}\sigma _z\mathrm{cos}\nu t`$.
Terms that modify absorption at the original linear resonance, $`\chi ^{\prime \prime }P_{n0}^{\prime \prime }(\nu )`$ can be grouped as follows:
$`\mathrm{\Omega }_\gamma =`$ $`{\displaystyle \frac{1}{4}}({\displaystyle \frac{h}{2}})^4\nu {\displaystyle \frac{1}{\mathrm{}^3}}{\displaystyle \underset{nml}{}}\sigma _{0n}\sigma _{nm}\sigma _{ml}\sigma _{l0}`$ (A17)
$`\times [P_{0n}^{\prime \prime }(\nu )P_{0m}^{}(2\nu )P_{0l}^{}(\nu )`$
$`+P_{0n}^{\prime \prime }(\nu )P_{0m}^{}(0)(P_{0l}^{}(\nu )P_{l0}^{}(\nu ))]`$
The two photon absorption terms can also be grouped as follows:
$`\mathrm{\Omega }_{2\gamma }=`$ $`{\displaystyle \frac{1}{4}}({\displaystyle \frac{h}{2}})^4\nu {\displaystyle \frac{1}{\mathrm{}^3}}{\displaystyle \underset{nml}{}}\sigma _{0n}\sigma _{nm}\sigma _{ml}\sigma _{l0}`$ (A19)
$`\times P_{0m}^{\prime \prime }(2\nu )P_{l0}^{}(\nu )P_{n0}^{}(\nu )`$ |
warning/0002/nlin0002006.html | ar5iv | text | # Turbulence Driven by a Deterministic Chaotic Dynamics
## Abstract
In the inertial range of fully developed turbulence, we model the vertex network dynamics by an iterated unimodular map having the universal behavior. Inertial range anomalous scaling for the pair correlation functions of the velocity and the local energy dissipation is established as a consequence of the chaotic behavior of the unimodular map when the Feigenbaum attractor looses stability. The anomalous corrections determined by the Feigenbaum constant $`\eta `$ to the Kolmogorov’s spectra are larger than those observed in experiments.
Turbulence Modeling and Turbulence Control in fluid dynamics is a paradigmatic problem and successful techniques in this domain may serve as models to deal with other complex extended systems. It is generally accepted now that a real understanding of turbulence requires calculations and models that are globally sensitive to all length scales.
To elucidate this conviction, let us briefly explain the present situation in the quantum field theory of fully developed turbulence of incompressible fluid which has been developed for the last two decades (see for a review). In order to be specific we consider the correlator $`K=v(𝐤)v(𝐤)`$ of the velocity field $`v`$ in the momentum representation ($`k|𝐤|`$ and $`\mathrm{}`$ denote the time averaging). It is found from the Dyson equation
$$K^1=\nu k^2+m^2\mathrm{\Sigma }\left(k\right),$$
(1)
where $`m`$ is the inverse integral turbulent scale, $`\nu `$ is the viscosity parameter, and $`\mathrm{\Sigma }(k)`$ is an infinite sum of all 1-irreducible diagrams of the Wyld’s diagrammatic technique . In the framework of $`42\epsilon `$ expansion the series (1) takes the form
$$K^1=\left(\nu k^2+m^2\right)\left[1+\underset{n=1}{\overset{\mathrm{}}{}}\left(gk^{2\epsilon }\right)^nc_n(m/k,\epsilon )\right],$$
(2)
The inverse viscous dissipative length $`\lambda `$ plays the role of a maximum momentum in the problem. The Reynolds number Re$`=(\lambda /m)^{4/3}`$ is assumed to be large. For $`km\lambda `$ and $`\epsilon >0`$, the dimensionless parameter of the expansion $`gk^{2\epsilon }(k/\lambda )^{2\epsilon }/\nu `$ in (2) is not small and it is necessary to sum the series. This problem reduces to a determination of the asymptotic value of $`K_\chi =K(\chi k,\chi m)`$ for $`\chi 0`$ (everything is fixed except for $`\chi `$). It is solved by the renormalization group procedure (RG) for any finite ratio $`m/k`$ -.
However, the coefficient $`c_n`$ in (2) has additional singularities as $`m/k0`$ which cannot be handled by RG since they have nothing to do with the scale invariance property (the parameter $`m`$ does not involve the renormalization). These singularities express the intermittency phenomenon (, Chap. 25) that is the actual reason why the bare RG is ineffective in turbulence .
The $`m/k0`$ asymptotics of the series (2) is found due to the Short Distance Expansion method , introduced by Wilson. It leads to the form
$$K\left(k\right)=Ak^{11/3}\left[1+\underset{i}{}b_i\left(m/k\right)^{\mathrm{\Delta }_{F_i}}\right]$$
(3)
where $`A`$ is the Kolmogorov’s constant, $`b_i`$ is the coefficient analytic in $`m/k`$, and $`\mathrm{\Delta }_{F_i}`$ is the critical exponent of a set of composite operators $`F_i`$. The series in (3) would give the corrections to the leading Kolmogorov’s asymptotics if all $`\mathrm{\Delta }_{F_i}>0`$.
A number of works has been devoted to the computation of $`\mathrm{\Delta }_{F_i}`$ for various sets of composite operators in hydrodynamics , and magnetohydrodynamics . It is found that there are infinitely many composite operators for which $`\mathrm{\Delta }_{F_i}0.`$ An example of such a “dangerous” one is the composite operator of local energy dissipation $`E=\nu (_iv_k+_kv_i)^2/2`$ with $`\mathrm{\Delta }_E=0`$ . The contributions to (3) due to such operators are not small and it is necessary to sum the infinite series in (3). There is no reliable technique for such a summation found. To proceed further, remaining of the idea of homogeneously fully developed turbulence, one must assign the statistics to each set of composite operators with $`\mathrm{\Delta }_{F_i}0.`$
It is conventional now that one has to employ a model, in which all scales are presented equally to achieve the further progress. In this Letter we suggest such a model.
The hydrodynamics of the incompressible fluid is described by the Navier-Stokes equation
$$_tv_i=\nu \mathrm{\Delta }v_i,_t_t+\left(v\right),$$
(4)
where $`v_i(𝐱,t)`$ is the transverse ($`_iv_i=0`$) velocity field, $`\nu `$ is the viscosity. To meet the energy balance in the system, one has to compensate the viscous dissipation by an energy pump of power $`W`$ with some spectral density $`𝔡(k)`$,
$$W=\frac{1}{\left(2\pi \right)^3}𝑑𝐤𝔡\left(k\right).$$
(5)
In - and in the descendant papers -, the pump source has been introduced in (4) by the Gaussian white noise $`f`$. The power model for the spectral density $`𝔡(k)k(k^2+m^2)^\epsilon `$ chosen in literature defines the Gaussian statistics for $`f`$ completely, $`ff(k)𝔡(k)`$. Nevertheless, such a model is not uniquely determined: One can introduce an arbitrary bounded value function $`𝔥(m/k)`$ ($`𝔥(0)=1`$) into the Gaussian covariance. This function did not involve the renormalization procedure in -. From the physical point of view, this fact means that the field $`f`$ has some redundant degrees of freedom which do not change to the scaling ones in the process of renormalization.
In actual experiments, the power $`W`$ pumped into the inertial range of scales $`mk\lambda `$ is delivered by the large scale vertices entering the system sometime in a variety of sizes ($`km`$), energies, and enstrophies. These vertices dissolve due to the strong nonlinearity of hydrodynamical interactions into ones of the smaller scales ($`k>m`$) forming the direct energy cascade. Alternatively, a small scale vertex would associate with others in the fractal sets contributing towards the inverse cascade. The resulting picture constitutes a vortex network encoded, at each moment of time, by a discrete set of ratios $`0<\{m/k\}_n<1`$. This network has much in common with dynamical systems.
We shall consider the pumping generated by some dynamical system,
$$𝒲_\tau \left(t\right)=\tau ^{1/2}\underset{n=0}{\overset{\left[t/\tau \right]}{}}z_n.$$
(6)
Here $`\tau >0`$ is a time step, the square brackets $`[\mathrm{}]`$ denote the integer part. $`z_n`$ are the iterates of a $`1`$parameter unimodular map $`T:𝕏𝕏`$ of the phase space $`𝕏[0,1]`$. We also suppose that the map $`T`$ possesses the invariant measure $`\mu _T.`$
As an example of such a map, one can consider the triangular map $`Tz=r\left(12\left|\frac{1}{2}z\right|\right)`$ or the logistic map $`Tz=rz(1z)`$ which is known to describe the angles of a strongly damped kicked rotator. However, the existence of invariant measure for the latter map is not proven for any $`r`$ except $`r=4`$. The dynamics of unimodular maps has been studied extensively last decades. For the values of the control parameter $`1<r<r_{\mathrm{}},`$ the relevant Liapunov exponent is always negative (except the bifurcation points $`r_n`$ when it becomes zero), while for $`r_{\mathrm{}}<r`$ this exponent is mostly positive indicating chaotic behavior.
The chaotic behavior observed in the unimodular maps results from the successive pitchfork bifurcations, which provides the mechanism for the successive doubling of the fixed points. It is important that the correlations between the $`j`$-shifted sequences
$$𝒞\left(j\right)=\underset{\tau 0}{lim}\left[\frac{\tau }{tt^{}}\right]\underset{n>0}{\overset{\left[tt^{}/\tau \right]}{}}z_nz_{n+j}$$
(7)
generated by the unimodular map decay with a power law in $`j`$ as $`rr_{\mathrm{}}`$ ,
$$𝒞\left(j\right)j^\eta 𝒞(1,\rho )+𝒪\left(rr_{\mathrm{}}\right).$$
(8)
Here $`\eta `$ is the universal constant , $`\rho `$ is the “scaling function” invariant with respect to the doubling transformation $`T^2`$.
An external force $`𝔣_\tau (𝐱,t):^{d+1}\times 𝕏^{d+1},`$ which we introduce in (4) is regarded as a priori arbitrary chaotic process $`\dot{𝒲}_\tau \tau ^{1/2}_{n0}z_n\delta (tn\tau ),`$ $`z_n𝕏,`$ generated by a deterministic chaotic dynamics of the map $`T`$ ($`d`$ is the space dimensionality).
In the language of quantum field theory, the time step $`\tau >0`$ plays the role of the limiting time scale, dividing the statistics of “fast” ($`t<\tau `$) and “slow” ($`t>\tau `$) modes. It is evident that the statistics of the fast modes $`𝔣(𝐱,t<\tau )`$ is not affected by the dynamics of the map $`T`$, so that one can treat as a Gaussian white noise with the covariance
$$D_{rs}(𝐱𝐲,tt^{})=\left(2\pi \right)^3𝑑\omega 𝑑𝐤P_{rs}𝔡\left(k\right)e^{i\left[𝐤\left(𝐱𝐲\right)\omega \left(tt^{}\right)\right]}$$
(9)
where $`P_{sr}=\delta _{sr}k_sk_r/k^2`$ is the transversal projector. Following , we chose the spectral density in the power form $`𝔡(k)=g\nu ^3k^{4d2\epsilon },`$ in which $`g`$ is a coupling constant. The actual value of the parameter of regular expansion is $`\epsilon =2,`$ so that $`𝔡(k)`$ represents a power model for $`\delta (𝐤)`$ located in the infrared ($`k0`$) region. The other way round, the statistics of the slow modes $`𝔣_\tau (𝐱,t>\tau )`$ is essentially non-Gaussian.
Then the velocity field correlator $`K=𝐯(𝐱,t)𝐯(𝐲,t^{})`$ factors in
$$K=𝔎(𝐱𝐲,tt^{})\underset{n2}{}𝒞_\tau {}_{n}{}^{}(tt^{};𝔮_1,\mathrm{}𝔮_n),t>t^{}.$$
(10)
Here $`𝔎(𝐱𝐲,tt^{})`$ is the correlator in the fast modes Gaussian statistics, and the infinite sum describes the chaotic process as $`t>\tau `$. All correlations
$$𝒞_\tau {}_{n}{}^{}=\underset{\tau 0}{lim}\left[\frac{\tau }{tt^{}}\right]\underset{𝔭>0}{\overset{\left[tt^{}/\tau \right]}{}}z_𝔭z_{𝔭+𝔮_1}\mathrm{}z_{𝔭+𝔮_n},𝔮_k\left[\frac{tt_k}{\tau }\right],$$
(11)
$`\left(t>\mathrm{}>t_k>t_{k1}>\mathrm{}t^{}>0\right)`$ exist since there is the continuous invariant measure $`\mu _T`$.
The correlator (10) is finite for any ratio $`m/k`$ fixed: $`𝔎`$ is found from the Dyson equation (1) in the framework of some quantum field theory with the multiplicatively renormalized action functional
$$S_𝐑=\frac{1}{2}𝔐^{2\epsilon }𝑑𝐤𝑑t𝐯^{}D𝐯^{}$$
(12)
$$+𝑑𝐤𝑑t𝐯^{}\left[_t𝐯\left(v\right)𝐯+\nu Z_\nu \mathrm{\Delta }𝐯\right]$$
of the basic $`𝐯`$ and auxiliary $`𝐯^{}`$ fields. $`D`$ is the covariance (9). The only renormalization constant $`\nu _0=\nu Z_\nu `$ is required to subtract singularities, $`𝔐`$ is the renormalization mass parameter. All correlation functions of the renormalized quantum field theory are finite for any fixed $`m/k`$. In the large-scale asymptotics $`sk/𝔐0,`$ in $`d=3`$, they demonstrate the scaling behavior with the Kolmogorov’s critical indices of velocity $`\mathrm{\Delta }_v=1/3`$ and time $`\mathrm{\Delta }_t=2/3.`$ As we have mentioned above, in the limit $`m/k0,`$ the additional singularities arise in $`𝔎`$. We demonstrate, however, that (10) is still finite.
In the majority of experimental data, the correlator $`K`$ and the correlation function of the local energy dissipation $`E\mathrm{\Phi }(𝐱)\mathrm{\Phi }(𝐲)`$, $`\mathrm{\Phi }(𝐱)=_iv_j(𝐱)_jv_i(𝐱),`$ have a power law behavior
$$K\left(k\right)k^{\left(11/3+\delta \right)},E(𝐱,𝐲)\left|𝐱𝐲\right|^\lambda $$
(13)
with small positive values of indices $`\lambda 0.2,`$ $`0.02<\delta <0.07`$ in accordance with ,. The Kolmogorov’s theory predicts $`\lambda =\delta =0.`$ Now we demonstrate that in the scaling limit $`\tau 0,`$ at $`r_{\mathrm{}}`$ (i.e., when the Feigenbaum attractor of the map $`T`$ looses its stability) the infinite sum in (10) provides the anomalous scaling $`\lambda >\delta >0.`$
First, we show that the pair correlation $`𝒞_2_\tau `$ provides the leading contribution to the sum in (10) $`(\tau /tt^{})^\eta `$. The contributions due to the triple and quadruple correlations etc. decay much more rapidly as $`\tau 0.`$ To prove it in an elegant way, we use the theorem followed which is analogous to the Wick’s theorem of quantum field theory: The average of the product is obtained as a sum of all possible paired averages.
Theorem
In the scaling region $`\tau 0,`$ the following expansion has place
$$𝒞_\tau {}_{n}{}^{}(t;𝔮_1,\mathrm{}𝔮_n)=\underset{i,j=1}{\overset{n}{}}𝒞_2{}_{\tau }{}^{}(t;𝔮_i)\underset{j>s,si}{}𝒞_2{}_{\tau }{}^{}(t;𝔮_j𝔮_s).$$
The elementary proof of the theorem is given in the Appendix. This theorem gives a key to diagram technique in dynamical systems. The diagrammatic approach to computation of arbitrary correlation functions in dynamical system theory has been developed recently in . One can estimate the scaling asymptotics for correlation functions as $`𝒞_\tau {}_{n}{}^{}\tau ^{n\eta /2}`$ ($`n`$ is assumed to be an even number) then the result on the scaling behavior seems obvious.
Interesting in the scaling asymptotics $`\tau 0`$, we neglect all higher order correlations in (10) except $`𝒞_2_\tau `$, so that, for $`rr_{\mathrm{}}`$, in the inertial range, in $`d=3`$, (10) reads as follows
$$K(k,tt^{})=A^{}k^{11/3}\left(\frac{\tau }{tt^{}}\right)^\eta +𝒪\left(rr_{\mathrm{}}\right),$$
(14)
in which $`A^{}`$ is some modified Kolmogorov’s constant. The asymptotics (14) can be treated as a consequence of modification of the critical regime in the theory (12) as the result of rescaling:
$$\stackrel{~}{v}=v\xi ,\stackrel{~}{v}^{}=v^{}\xi ^3,\stackrel{~}{g}=g\xi ^{2+d},\stackrel{~}{k}=k\xi ,\stackrel{~}{t}=t\xi ^2,$$
(15)
in which
$$\xi \left(\tau /t\right)^\sigma ,\sigma =\frac{\eta }{2+d2\mathrm{\Delta }_v},$$
$`\mathrm{\Delta }_v`$ is the Kolmogorov’s dimension of velocity. The critical exponents of all $`\stackrel{~}{}`$-quantities are known (they equal to the Kolmogorov’s values). The scaling parameter of the doubling transformation $`\tau `$ does not involve the theory (12), so that it has no definite critical dimension. Therefore, the modification of scaling asymptotics $`sk/𝔐0`$ due to (15) goes from the rescaling with respect to the dimensional parameter $`ts^{2/3}`$. The upper critical dimension (above which the hydrodynamical interaction becomes infrared unimportant, and the field theory is asymptotically free) for the rescaled model is changed to $`44/3(2+d)\sigma `$ and is not equal to the logarithmic dimension $`d_l=4`$. In dimension $`d=42\epsilon `$, the hydrodynamical interaction is important for $`2/3(2+d)\sigma <\epsilon `$ (this condition is satisfied for the actual value $`\epsilon =2`$). Thus, the quantities
$$\mathrm{\Delta }_k^{}=12/3\sigma ,\mathrm{\Delta }_v^{}=\mathrm{\Delta }_v2/3\sigma ,\mathrm{\Delta }_t^{}=\mathrm{\Delta }_t+4/3\sigma ,$$
(16)
generated by (15) are the entire critical dimensions of momentum, velocity, and time in turbulence launched by the dynamical system.
The simple power counting in (14) gives the asymptotics
$$Kk^{\left(11/3+\delta _\eta \right)},\delta _\eta =2\eta /3\sigma 0.2014722\mathrm{}.$$
The renormalization of the local energy dissipation operator $`E`$ with mixing has been given in in details. It was shown that the marginal contribution into asymptotics is risen due to the linear combinations of composite operators including $`𝐯\mathrm{\Delta }𝐯(𝐱),`$ which has the critical exponent $`\mathrm{\Delta }_E=43\gamma _t=0,`$ where $`\gamma _t=2+\mathrm{\Delta }_t=4/3`$ is the anomalous dimension of time in the Kolmogorov’s theory. In the rescaled theory, the anomalous dimension of time is changed to $`\gamma _t^{}=2+\mathrm{\Delta }_t^{}0.7107451\mathrm{},`$ so that the resulting critical dimension $`\mathrm{\Delta }_E^{}1.8677646\mathrm{}>0`$ is no more marginal. This corresponds to the power law behavior
$$E(𝐱,𝐲)\left|𝐱𝐲\right|^\lambda ,\lambda 1.1322354\mathrm{}.$$
In conclusion, we have studied the model of turbulence driven by the deterministic chaotic dynamics generated by a unimodular iterated map. The critical scaling is established when the Feigenbaum attractor looses stability. When the new time scale comes into play, the Kolmogorov’s critical regime is modified and the anomalous exponents becomes positive $`\lambda >\delta >0.`$ However, their absolute values prescribed by the Feigenbaum constant $`\eta `$ are larger than those in ,.
Appendix
A proof sketch of the Theorem.
Let us consider the product of partial sums
$$𝔓(𝔮,𝔮_1\mathrm{}𝔮_n)$$
$$\frac{1}{n}\underset{𝔭}{}z_𝔭z_{𝔭+𝔮}\frac{1}{n}\underset{𝔭_1}{}z_{𝔭_1+𝔮_1}z_{𝔭_1+𝔮_2}\mathrm{}\frac{1}{n}\underset{𝔭_n}{}z_{𝔭_n+𝔮_{n1}}z_{𝔭_n+𝔮_n}.$$
Rename the counting variable in each sum, for example
$$\frac{1}{n}\underset{𝔭_1}{}z_{𝔭_1+𝔮_1}z_{𝔭_1+𝔮_2}=\frac{n+𝔮_1}{n}\frac{1}{n+𝔮_1}\underset{𝔯=1+𝔮_1}{\overset{n+𝔮_1}{}}z_𝔯z_{𝔯+\left(𝔮_2𝔮_1\right)}.$$
As $`n\mathrm{}`$, this partial sum tends to $`𝒞_2(𝔮_2𝔮_1).`$ It remains to note that
$$𝒞_n(𝔮_1,\mathrm{}𝔮_n)=\underset{𝔮_1\mathrm{}𝔮_{n1}}{}\delta _{𝔭𝔭_1}\mathrm{}\delta _{𝔭𝔭_{n1}}𝔓(𝔮,𝔮_1\mathrm{}𝔮_n).$$ |
warning/0002/gr-qc0002047.html | ar5iv | text | # Peeling Properties of Light–Like Signals in General Relativity
## 1 Introduction
The history of a light–like signal in General Relativity is a singular null hypersurface. The null hypersurface is called singular because in general the Ricci tensor and the Weyl tensor of the space–time contain Dirac $`\delta `$ function terms with the $`\delta `$function singular on the null hypersurface. Such a singular null hypersurface can be used as a simplified model of a supernova if the space–time before and after the emission of the light–like signal is a model of the vacuum field due to an isolated gravitating system and if the singular null hypersurface is asymptotically (as future null infinity is approached) a future–directed null–cone. For the model described in the space–times before and after the emission of the light–like signal are two copies of the Weyl asymptotically flat, static space–times but with different multipole moments. Thus the explosion is modelled by a sudden change in the multipole moments of the source. The coefficients of the $`\delta `$ function in the Weyl tensor display unconventional peeling properties. By peeling properties here we mean the dependence of the coefficients of the $`\delta `$function in the Weyl tensor on an affine parameter $`r`$ along the generators of the singular null hypersurface, with $`r+\mathrm{}`$ as future null infinity is approached. In general the $`\delta `$ function in the Weyl tensor can be unambiguously split into a matter part (if it exists) of Petrov type II and a wave part (if it exists) of Petrov type N. This splitting, which was originally announced in , is fully developed in . In the simplified supernova model described in the matter part of the $`\delta `$function in the Weyl tensor depends on $`r`$ in the form of $`O(r^3)`$– terms and smaller terms, while the wave part of the $`\delta `$ function in the Weyl tensor has a coefficient which is $`O(r^4)`$. We denote the components of the $`\delta `$–function in the Weyl tensor, in Newman–Penrose notation, by $`\widehat{\mathrm{\Psi }}_A(A=0,1,2,3,4)`$ chosen in such a way that the first four of these describe the matter part of the Weyl tensor and $`\widehat{\mathrm{\Psi }}_4`$ describes the wave part. Defining $`\widehat{\mathrm{\Psi }}_A`$ in this way the conventional peeling behavior would be to have $`\widehat{\mathrm{\Psi }}_A=O(r^{5+A})(A=0,1,2,3,4)`$. This is the conventional peeling behavior when compared with models of vacuum gravitational fields due to isolated gravitating systems in general (cf. and (2.20)–(2.24) below). The purpose of the present paper is to explain the unconventional peeling behavior described above by putting it in the context of light–like signals emitted by a general class of isolated gravitating systems.
We take the vacuum gravitational field outside an isolated system, before and after the emission of a light–like signal, to be modelled by a Bondi–Sachs asymptotically flat space–time. We also consider the most general matching of the two space–times on the singular null hypersurface separating them. How the two space–times are glued together influences the type of signal whose history is the singular null hypersurface. We demonstrate that in general the peeling behavior is the conventional one which is associated with a radiating isolated system and that it becomes unconventional if the asymptotically flat space–times on either side of the history of the light–like signal tend to flatness at future null infinity faster than the general Bondi–Sachs space–time.
## 2 Light–Like Signal From Isolated Source
Throughout this paper the space–time model of the vacuum gravitational field outside an isolated source will be a Bondi–Sachs , space–time. A convenient form of this space–time line–element is given by
$$ds^2=(\theta ^1)^2+(\theta ^2)^22\theta ^3\theta ^4,$$
(2.1)
with
$`\theta ^1`$ $`=`$ $`rp^1(\mathrm{e}^\alpha \mathrm{cosh}\beta dx+\mathrm{e}^\alpha \mathrm{sinh}\beta dy+adu),`$ (2.2)
$`\theta ^2`$ $`=`$ $`rp^1(\mathrm{e}^\alpha \mathrm{sinh}\beta dx+\mathrm{e}^\alpha \mathrm{cosh}\beta dy+bdu),`$ (2.3)
$`\theta ^3`$ $`=`$ $`dr{\displaystyle \frac{1}{2}}cdu,`$ (2.4)
$`\theta ^4`$ $`=`$ $`du.`$ (2.5)
The six function $`\alpha ,\beta ,a,b,p,c`$ depend on all coordinates $`x,y,r,u`$ and $`u=\mathrm{const}.`$ are null hypersurfaces generated by the geodesic integral curves of the vector field $`\frac{}{r}`$ with $`r`$ an affine parameter along these geodesics. The following assumptions are made regarding the $`r`$–dependence of these functions:
$`\alpha `$ $`=`$ $`{\displaystyle \frac{\alpha _1}{r}}+{\displaystyle \frac{\alpha _2}{r^2}}+{\displaystyle \frac{\alpha _3}{r^3}}+\mathrm{},`$ (2.6)
$`\beta `$ $`=`$ $`{\displaystyle \frac{\beta _1}{r}}+{\displaystyle \frac{\beta _2}{r^2}}+{\displaystyle \frac{\beta _3}{r^3}}+\mathrm{},`$ (2.7)
$`a`$ $`=`$ $`a_0+{\displaystyle \frac{a_1}{r}}+{\displaystyle \frac{a_2}{r^2}}+{\displaystyle \frac{a_3}{r^3}}+\mathrm{},`$ (2.8)
$`b`$ $`=`$ $`b_0+{\displaystyle \frac{b_1}{r}}+{\displaystyle \frac{b_2}{r^2}}+{\displaystyle \frac{b_3}{r^3}}+\mathrm{},`$ (2.9)
$`p`$ $`=`$ $`p_0\left(1+{\displaystyle \frac{q_1}{r}}+{\displaystyle \frac{q_2}{r^2}}+{\displaystyle \frac{q_3}{r^3}}+\mathrm{}\right),`$ (2.10)
$`c`$ $`=`$ $`1{\displaystyle \frac{2m}{r}}+\mathrm{}.`$ (2.11)
Here $`p_0=1+\frac{1}{4}(x^2+y^2)`$ and all the other coefficients of the inverse powers of $`r`$ displayed above are functions of $`(x,y,u)`$. All eighteen functions of $`(x,y,u)`$ appearing in (2.6)–(2.11) are required in order to have a knowledge of the metric tensor components up to and including $`\frac{1}{r}`$–terms. The vacuum field equations and an outgoing radiation condition allow us to specialise the sixteen functions as follows:
$$\alpha _2=\beta _2=0,$$
(2.12)
$$a_0=a_1=0,\mathrm{and}b_0=b_1=0,$$
(2.13)
$`a_2`$ $`=`$ $`p_0^4\left\{{\displaystyle \frac{}{x}}(p_0^2\alpha _1)+{\displaystyle \frac{}{y}}(p_0^2\beta _1)\right\},`$ (2.14)
$`b_2`$ $`=`$ $`p_0^4\left\{{\displaystyle \frac{}{x}}(p_0^2\beta _1){\displaystyle \frac{}{y}}(p_0^2\alpha _1)\right\},`$ (2.15)
$$q_1=0,q_2=\frac{1}{2}(\alpha _1^2+\beta _1^2),q_3=0.$$
(2.16)
When these equations are satisfied there remain five further field equations. They are propagation equations for $`m(u,x,y),\alpha _3(u,x,y),\beta _3(u,x,y),a_3(u,x,y)`$ and $`b_3(u,x,y)`$ off $`u=\mathrm{const}`$. The simplest reads
$$\dot{M}+\left|\dot{\gamma }\right|^2=0,$$
(2.17)
where the dot denotes partial differentiation with respect to $`u`$, and
$$M=m\dot{q}_2\frac{1}{2}\left\{\frac{}{x}(p_0^2a_2)+\frac{}{y}(p_0^2b_2)\right\},$$
(2.18)
with
$$\gamma =\alpha _1+i\beta _1.$$
(2.19)
When the equation (2.17) is averaged over the 2–sphere with line–element $`dl^2=p_0^2(dx^2+dy^2)`$ the well–known Bondi–Sachs mass–loss formula results. The remaining equations involving $`\dot{\alpha }_3,\dot{\beta }_3,\dot{a}_3`$ and $`\dot{b}_3`$ are given in and will not be used here. The curvature tensor components, in Newman– Penrose notation, for the space–time described above display the conventional peeling behavior:
$`\mathrm{\Psi }_0`$ $`=`$ $`{\displaystyle \frac{1}{r^5}}\left\{6(\alpha _3+i\beta _3){\displaystyle \frac{3}{2}}(\gamma +\overline{\gamma })^2(\gamma \overline{\gamma })2\overline{\gamma }^3\right\}+\mathrm{},`$ (2.20)
$`\mathrm{\Psi }_1`$ $`=`$ $`{\displaystyle \frac{1}{r^4\sqrt{2}}}\left\{{\displaystyle \frac{3}{2}}p_0^1(a_3+ib_3)+3p_0^3\gamma {\displaystyle \frac{}{\overline{z}}}(p_0^2\overline{\gamma })\right\}+\mathrm{},`$ (2.21)
$`\mathrm{\Psi }_2`$ $`=`$ $`{\displaystyle \frac{1}{r^3}}\left\{M+\gamma {\displaystyle \frac{\overline{\gamma }}{u}}+2p_0^2{\displaystyle \frac{}{\overline{z}}}\left({\displaystyle \frac{}{\overline{z}}}(p_0^2\overline{\gamma })\right)\right\}+\mathrm{},`$ (2.22)
$`\mathrm{\Psi }_3`$ $`=`$ $`{\displaystyle \frac{2}{r^2\sqrt{2}}}p_0^2{\displaystyle \frac{}{u}}\left(p_0{\displaystyle \frac{}{\overline{z}}}(p_0^2\overline{\gamma })\right)+\mathrm{},`$ (2.23)
$`\mathrm{\Psi }_4`$ $`=`$ $`{\displaystyle \frac{1}{r}}{\displaystyle \frac{^2\overline{\gamma }}{u^2}}+\mathrm{},`$ (2.24)
where we have put $`z=x+iy`$ and a bar denotes complex conjugation. Finally the complex shear $`\sigma `$ and the real expansion $`\vartheta `$ of the null geodesic integral curves of the vector field $`\frac{}{r}`$ are given by
$`\sigma `$ $`=`$ $`{\displaystyle \frac{(\alpha _1+i\beta _1)}{r^2}}{\displaystyle \frac{3(\alpha _3+i\beta _3)+2\alpha _1\beta _1^2}{r^4}}+O\left(r^5\right),`$ (2.25)
$`\vartheta `$ $`=`$ $`{\displaystyle \frac{1}{r}}+{\displaystyle \frac{2q_2}{r^3}}+O\left(r^5\right),`$ (2.26)
respectively, demonstrating that asymptotically (as $`r+\mathrm{}`$) the null hypersurfaces $`u=\mathrm{const}.`$ are future–directed null–cones.
We now consider the space–time above to be subdivided into two halves $`M^{}(u0)`$ and $`M^+(u0)`$ each with boundary the null hypersurface $`u=0`$. Let $`x_+^\mu =(x_+,y_+,r_+,u)`$ (with $`\mu =1,2,3,4`$) be the local coordinate system in $`M^+`$ in which the line–element takes the form given by (2.1)–(2.5) with the coefficients of the powers of $`r_+^1`$ in (2.6)–(2.11) denoted with a superscript plus as $`\alpha _1^+,\alpha _2^+`$, etc.. Let $`x^\mu =(x,y,r,u)`$ be the local coordinate system in $`M^{}`$ in terms of which the line–element of the space–time has the form given by (2.1) –(2.5). We take $`\xi ^a=(x,y,r)`$, with $`a=1,2,3`$, as local intrinsic coordinates on $`u=0`$. Now $`u=0`$ is the history of a light–like signal emitted by the isolated source and propagating into the space–time $`M^{}`$ leaving the space–time $`M^+`$ behind. We apply the Barrabès–Israel technique (see also for some recent developments) to analyse the physical properties of the signal with history $`u=0`$. We shall assume that the reader is familiar with . While we are applying the technique to the current situation we will comment on what we are doing so as to guide the reader by example through . The first requirement of is that the metric tensors induced on the null hypersurface $`u=0`$ by its embedding in $`M^+`$ and in $`M^{}`$ agree on $`u=0`$. This is achieved if the space–times $`M^+`$ and $`M^{}`$ described above are attached on $`u=0`$ with the following matching conditions:
$`x_+`$ $`=`$ $`f(x,y)+{\displaystyle \frac{f_1(x,y)}{r}}+O\left(r^2\right),`$ (2.27)
$`y_+`$ $`=`$ $`g(x,y)+{\displaystyle \frac{g_1(x,y)}{r}}+O\left(r^2\right),`$ (2.28)
$`r_+`$ $`=`$ $`rh(x,y)+h_0(x,y)+O\left(r^1\right),`$ (2.29)
with
$$f_x=g_y,f_y=g_x,h=\frac{P}{p_0}\left\{f_x^2+g_x^2\right\}^{\frac{1}{2}},$$
(2.30)
where $`P=1+\frac{1}{4}(f^2+g^2)`$ and the subscripts on $`f,g`$ denote partial differentiation. In addition we must have
$`hP^2(f_x^2f_y^2)\beta _1^++2hP^2f_xf_y\alpha _1^+p_0^2\beta _1=`$ (2.31)
$`{\displaystyle \frac{1}{2}}h^2P^2\left\{f_x\left({\displaystyle \frac{g_1}{x}}+{\displaystyle \frac{f_1}{y}}\right)+f_y\left({\displaystyle \frac{f_1}{x}}{\displaystyle \frac{g_1}{y}}\right)\right\},`$
$`hP^2(f_x^2f_y^2)\alpha _1^+2hP^2f_xf_y\beta _1^+p_0^2\alpha _1=`$ (2.32)
$`{\displaystyle \frac{1}{2}}h^2P^2\left\{f_x\left({\displaystyle \frac{f_1}{x}}{\displaystyle \frac{g_1}{y}}\right)f_y\left({\displaystyle \frac{g_1}{x}}+{\displaystyle \frac{f_1}{y}}\right)\right\},`$
and
$`h_0+{\displaystyle \frac{1}{2}}p_0^1(f_x^2+f_y^2)^{\frac{1}{2}}(ff_1+gg_1)=`$ (2.33)
$`{\displaystyle \frac{1}{2}}Pp_0^1(f_x^2+f_y^2)^{\frac{3}{2}}\left\{f_x\left({\displaystyle \frac{f_1}{x}}+{\displaystyle \frac{g_1}{y}}\right)+f_y\left({\displaystyle \frac{f_1}{y}}{\displaystyle \frac{g_1}{x}}\right)\right\}.`$
Here $`\beta _1^+=\beta _1^+(f,g,0),\beta _1=\beta _1(x,y,0),\alpha _1^+=\alpha _1^+(f,g,0)`$ and $`\alpha _1=\alpha _1(x,y,0)`$. The complexity of these matching conditions suggests that we examine the light–like signal in two stages. The leading terms in (2.27)–(2.28) are constructed from the analytic function $`F(z)=f(x,y)+ig(x,y)`$ with $`z=x+iy`$. They describe a part of the gluing of $`M^+`$ to $`M^{}`$ on $`u=0`$ which is a Penrose warp . This particular gluing leads to an impulsive gravitational wave with history $`u=0`$ having a line or directional singularity and we will consider it in section 3. Thus for the remainder of this section we shall take $`f(x,y)=x`$ and $`g(x,y)=y`$ in (2.27) and (2.28) and thus in (2.29) $`h=1`$. Now (2.31)–(2.33) simplify to read
$$\left[\beta _1\right]=\frac{1}{2}\left(\frac{g_1}{x}+\frac{f_1}{y}\right),\left[\alpha _1\right]=\frac{1}{2}\left(\frac{f_1}{x}\frac{g_1}{y}\right),$$
(2.34)
and
$$h_0=\frac{1}{2}p_0^2\frac{}{x}(p_0^2f_1)\frac{1}{2}p_0^2\frac{}{y}(p_0^2g_1),$$
(2.35)
where the square brackets will henceforth denote the jump across $`u=0`$ of the quantity within them, calculated in the coordinates $`x_{}^\mu `$. Thus, for example, $`\left[\beta _1\right]=\beta _1^+(x,y,u=0)\beta _1(x,y,u=0)`$.
The Barrabès–Israel technique is an extension to hypersurfaces of all types of the extrinsic curvature technique for non–null hypersurfaces (see, for example ). In the case of a null hypersurface the normal is tangent to the hypersurface so in order to obtain an analogous quantity to extrinsic curvature one first constructs a ‘transverse vector field’ which is any vector field defined on the hypersurface which is not tangent to the hypersurface and which is the same vector field when viewed in the coordinates $`(x_+^\mu )`$ and in the coordinates $`(x_{}^\mu )`$. In our case the normal to $`u=0`$ is given via the 1–form
$$n_\mu dx^\mu |_\pm =du,$$
(2.36)
where $`|_\pm `$ means the quantity is calculated in the plus or minus coordinates. A natural choice of transversal on the minus side, in view of the form of the line–element given by (2.1)–(2.5), is
$${}_{}{}^{}N_{\mu }^{}dx_{}^\mu =dr\frac{1}{2}cdu,$$
(2.37)
with $`c`$ given by (2.11). To ensure that the transversal when viewed on the plus side, $`{}_{}{}^{+}N_{\mu }^{}`$, is the same covariant vector field as $`{}_{}{}^{}N_{\mu }^{}`$ we proceed as follows: We pointed out above that we may use $`\xi ^a=(x,y,r)`$ as intrinsic coordinates on $`u=0`$. We then have three linearly independent tangent vectors to $`u=0`$ given by $`\frac{}{\xi ^a}`$. On the minus side of $`u=0`$ these have components
$$e_{(a)}^\mu |_{}=\frac{x_{}^\mu }{\xi ^a}=\delta _a^\mu ,$$
(2.38)
and on the plus side their components are
$$e_{(a)}^\mu |_+=\frac{x_+^\mu }{\xi ^a},$$
(2.39)
with $`x_+^\mu `$ given in terms of $`\xi ^a`$ by (2.27)–(2.29) \[now with $`f=x,g=y,h=1`$\]. Now $`{}_{}{}^{+}N_{\mu }^{}`$ is chosen so that
$$\left[N_\mu e_{(a)}^\mu \right]=0=\left[N_\mu N^\mu \right].$$
(2.40)
With $`{}_{}{}^{+}N_{\mu }^{}`$ thus calculated the ‘transverse extrinsic curvature’ on the plus or minus sides of $`u=0`$ is defined by
$$𝒦_{ab}^\pm ={}_{}{}^{\pm }N_{\mu }^{}\left(\frac{e_{(a)}^\mu |_\pm }{\xi ^b}+{}_{}{}^{\pm }\mathrm{\Gamma }_{\lambda \sigma }^{\mu }e_{(a)}^\lambda |_\pm e_{(b)}^\sigma |_\pm \right)=𝒦_{ba}^\pm ,$$
(2.41)
where $`{}_{}{}^{\pm }\mathrm{\Gamma }_{\lambda \sigma }^{\mu }`$ are the components of the Riemannian connection calculated on the plus or minus sides of $`u=0`$. We define
$$\gamma _{ab}=2\left[𝒦_{ab}\right],$$
(2.42)
and this is independent of the choice of transversal . Now $`\gamma _{ab}`$ is extended to a 4–tensor $`\gamma _{\mu \nu }`$ on $`u=0`$ by padding out with zeros (the only requirement on the extension being that its projection tangential to $`u=0`$ be $`\gamma _{ab}`$). With
$$\gamma ^\mu =\gamma ^{\mu \nu }n_\nu ,\gamma ^{}=\gamma ^\mu n_\mu ,\gamma =g^{\mu \nu }\gamma _{\mu \nu },$$
(2.43)
calculated in the plus or minus coordinates (we leave out the designation $`|_\pm `$ in such situations), the coefficient of the delta function $`\delta (u)`$ in the Einstein tensor of the space–time $`M^+M^{}`$ gives the surface stress–energy tensor
$$16\pi \eta ^1S^{\mu \nu }=2\gamma ^{(\mu }n^{\nu )}\gamma n^\mu n^\nu \gamma ^{}g^{\mu \nu },$$
(2.44)
where $`\eta ^1=n^\mu N_\mu `$. The coefficient of the delta function $`\delta (u)`$ in the Weyl tensor components is then
$$\widehat{C}^{\kappa \lambda }{}_{\mu \nu }{}^{}=2\eta n^{[\kappa }\gamma ^{\lambda ]}{}_{[\mu }{}^{}n_{\nu ]}^{}16\pi \delta _{[\mu }^{[\kappa }S_{\nu ]}^{\lambda ]}+\frac{8\pi }{3}S_\alpha ^\alpha \delta _{\mu \nu }^{\kappa \lambda }.$$
(2.45)
If $`m^\mu `$ is a unit complex vector field defined on $`u=0`$ which is tangential to $`u=0`$ and also orthogonal to the transversal then the Newman–Penrose components of $`\widehat{C}^{\kappa \lambda }_{\mu \nu }`$ are given by
$`\widehat{\mathrm{\Psi }}_0=0,\widehat{\mathrm{\Psi }}_1=0,\widehat{\mathrm{\Psi }}_2={\displaystyle \frac{1}{6}}\eta \gamma ^{},`$ (2.46)
$`\widehat{\mathrm{\Psi }}_3={\displaystyle \frac{1}{2}}\eta \gamma _\mu \overline{m}^\mu ,\widehat{\mathrm{\Psi }}_4={\displaystyle \frac{1}{2}}\eta \gamma _{\mu \nu }\overline{m}^\mu \overline{m}^\nu .`$
In general the signal is Petrov type II and contains a gravitational wave if $`\widehat{\mathrm{\Psi }}_40`$.
Carrying out this procedure (the calculations in this paper have been performed using GRTensorM version 1.2 for MATHEMATICA 3.x ) we find that $`u=0`$ has a non–vanishing surface stress–energy tensor with components
$`S^{11}`$ $`=`$ $`O\left({\displaystyle \frac{1}{r^5}}\right),`$ (2.47)
$`S^{22}`$ $`=`$ $`O\left({\displaystyle \frac{1}{r^5}}\right),`$ (2.48)
$`S^{12}`$ $`=`$ $`O\left({\displaystyle \frac{1}{r^6}}\right),`$ (2.49)
$`S^{13}`$ $`=`$ $`{\displaystyle \frac{1}{16\pi r^3}}\left\{2[a_2]+f_12p_0^2{\displaystyle \frac{h_0}{x}}2f_1\dot{\alpha }_1^+2g_1\dot{\beta }_1^+\right\}+\mathrm{},`$ (2.50)
$`S^{23}`$ $`=`$ $`{\displaystyle \frac{1}{16\pi r^3}}\left\{2[b_2]+g_12p_0^2{\displaystyle \frac{h_0}{y}}2f_1\dot{\beta }_1^++2g_1\dot{\alpha }_1^+\right\}+\mathrm{},`$ (2.51)
and
$$S^{33}=\frac{1}{4\pi r^2}\left\{[m\dot{q}_2]\frac{1}{2}p_0^2\frac{}{y}(p_0^2[b_2])\frac{1}{2}p_0^2\frac{}{x}(p_0^2[a_2])+\frac{1}{2}(h_0+\mathrm{\Delta }h_0)\right\}+\mathrm{},$$
(2.52)
where $`\mathrm{\Delta }=p_0^2\left(\frac{^2}{x^2}+\frac{^2}{y^2}\right)`$ and $`q_2`$ is given by (2.16). The jumps $`[a_2],[b_2]`$ can be written in terms of the jumps $`[\alpha _1],[\beta _1]`$ using the field equations (2.14) and (2.15) and thence in terms of the functions $`f_1,g_1`$ via (2.34) to arrive at
$`\left[a_2\right]`$ $`=`$ $`p_0{\displaystyle \frac{p_0}{x}}\left({\displaystyle \frac{g_1}{y}}{\displaystyle \frac{f_1}{x}}\right)+p_0{\displaystyle \frac{p_0}{y}}\left({\displaystyle \frac{g_1}{x}}+{\displaystyle \frac{f_1}{y}}\right){\displaystyle \frac{1}{2}}\mathrm{\Delta }f_1,`$ (2.53)
$`\left[b_2\right]`$ $`=`$ $`p_0{\displaystyle \frac{p_0}{x}}\left({\displaystyle \frac{g_1}{x}}+{\displaystyle \frac{f_1}{y}}\right)+p_0{\displaystyle \frac{p_0}{y}}\left({\displaystyle \frac{g_1}{y}}{\displaystyle \frac{f_1}{x}}\right){\displaystyle \frac{1}{2}}\mathrm{\Delta }g_1.`$ (2.54)
and finally at
$$S^{33}=\frac{1}{4\pi r^2}\left\{[m\dot{q}_2]+\frac{1}{4}p_0^2\left(\frac{}{x}(p_0^2f_1)+\frac{}{y}(p_0^2g_1)\right)\right\}+\mathrm{}.$$
(2.55)
The surface energy–density of the shell measured by a radially moving observer (see ) is a positive multiple of $`S^{33}`$. If we make the assumption that the functions $`p_0^2f_1,p_0^2g_1`$, defined on the 2–sphere with line–element $`dl^2=p_0^2(dx^2+dy^2)`$, are bounded then it follows from (2.55) that the leading term in the surface energy density averaged over the 2–sphere is proportional to the jump in the Bondi–Sachs mass across $`u=0`$ (this latter, as mentioned following (2.19), is the average of $`M`$ in (2.18) over the 2–sphere). The average surface energy–density is positive if there is a loss of Bondi–Sachs mass.
Finally we calculate $`\widehat{\mathrm{\Psi }}_A`$ for $`A=2,3,4`$ in (2.46). The result can be put in the form
$`\widehat{\mathrm{\Psi }}_2`$ $`=`$ $`O\left({\displaystyle \frac{1}{r^3}}\right),`$ (2.56)
$`\widehat{\mathrm{\Psi }}_3`$ $`=`$ $`4\pi \sqrt{2}rp_0^1\left(S^{13}iS^{23}\right)+O\left({\displaystyle \frac{1}{r^3}}\right)=O\left({\displaystyle \frac{1}{r^2}}\right),`$ (2.57)
and
$$\widehat{\mathrm{\Psi }}_4=\frac{[\dot{\alpha }_1i\dot{\beta }_1]}{r}+\frac{W}{r^2}+O\left(\frac{1}{r^3}\right),$$
(2.58)
with
$`W=2{\displaystyle \frac{}{z}}\left(p_0^2{\displaystyle \frac{h_0}{z}}\right){\displaystyle \frac{1}{2}}[\alpha _1i\beta _1]+{\displaystyle \frac{}{z}}[a_2ib_2]`$ (2.59)
$`+(\dot{\alpha }_1^+i\dot{\beta }_1^+)\left\{h_0i\left({\displaystyle \frac{g_1}{x}}{\displaystyle \frac{f_1}{y}}\right)\right\}.`$
From this we arrive at the main result of this section, namely, in general the coefficients of the delta function terms in the Weyl tensor display the conventional peeling behavior. The signal contains an impulsive gravitational wave part $`\widehat{\mathrm{\Psi }}_4`$ with the expected $`\frac{1}{r}`$–behavior provided $`[\dot{\alpha }_1i\dot{\beta }_1]0`$. This latter means that there is a jump in the Bondi ‘news’ across $`u=0`$. This impulsive wave is accompanied by a light–like shell with surface stress–energy tensor given by (2.47)–(2.55). We see from (2.58) and (2.59) that if the matching (2.27)–(2.29) is the identity matching then the wave produced by the jump in the news across $`u=0`$ is free from line singularities.
## 3 A General and a Special Example
We return now to the general matching conditions (2.27)–(2.33). A considerable computational effort is needed to establish the following orders of magnitude of the components of the stress–energy tensor on $`u=0`$ in this case:
$$S^{11}=O\left(\frac{1}{r^5}\right),S^{22}=O\left(\frac{1}{r^5}\right),S^{12}=O\left(\frac{1}{r^6}\right),$$
(3.1)
$$S^{13}=O\left(\frac{1}{r^3}\right),S^{23}=O\left(\frac{1}{r^3}\right),S^{33}=O\left(\frac{1}{r^2}\right).$$
(3.2)
In addition $`\widehat{\mathrm{\Psi }}_A`$ for $`A=2,3,4`$ in this case satisfy
$$\widehat{\mathrm{\Psi }}_2=O\left(\frac{1}{r^3}\right),\widehat{\mathrm{\Psi }}_3=O\left(\frac{1}{r^2}\right),$$
(3.3)
and
$$\widehat{\mathrm{\Psi }}_4=\frac{1}{r}\left\{p_0^2\left(H(z)\frac{(F^{})^2}{1+\frac{1}{4}|F(z)|^2}(\dot{\alpha }_1^+i\dot{\beta }_1^+)\right)+\dot{\alpha }_1i\dot{\beta }_1\right\}+O\left(\frac{1}{r^2}\right).$$
(3.4)
Here, as in the paragraph following (2.33), $`F(z)=f(x,y)+ig(x,y)`$ with $`z=x+iy`$, $`F^{}=dF/dz`$ and
$$H(z)=\frac{F^{\prime \prime \prime }}{F^{}}\frac{3}{2}\left(\frac{F^{\prime \prime }}{F}\right)^2.$$
(3.5)
Clearly when $`F(z)=z`$ this reduces to (2.58). We see from (3.3) and (3.4) that in this general case the conventional peeling behavior is exhibited. However we also see from the first $`\frac{1}{r}`$– term in (3.4) that the Penrose wave has a directional or line singularity (as $`z\overline{z}+\mathrm{}`$) .
We mentioned in the introduction that the principal motivation for the present study is to put into perspective the unconventional peeling behavior of $`\widehat{\mathrm{\Psi }}_A`$ for $`A=2,3,4`$ encountered in a simple example of a supernova in which the space–times $`M^+`$ and $`M^{}`$ are two copies of the Weyl asymptotically flat, static space– times having different multipole moments. We shall now demonstrate how this example emerges as a special case of the general situation described in section 2.
An asymptotically flat Weyl static space–time has a line–element which can be put in the form (2.1)–(2.11) with
$$\alpha _1=\beta _1=0,a_2=b_2=0,$$
(3.6)
$$\alpha _3=\frac{1}{2}Qp_0^2(x^2y^2),\beta _3=Qp_0^2xy,$$
(3.7)
$$a_3=2Dx,b_3=2Dy,$$
(3.8)
$`c=1{\displaystyle \frac{2m}{r}}2{\displaystyle \frac{D}{r^2}}p_0^2(1{\displaystyle \frac{1}{4}}(x^2+y^2))`$ (3.9)
$`{\displaystyle \frac{Q}{r^3}}p_0^2\left\{22(x^2+y^2)+{\displaystyle \frac{1}{8}}(x^2+y^2)^2\right\}+\mathrm{}.`$
The constant $`m`$ is interpreted as the mass of the source while the constants $`D`$ and $`Q`$ are taken to be the dipole and quadrupole moments of the source respectively. In $`M^+`$ the local coordinates are $`x_+^\mu =(x_+,y_+,r_+,u)`$ and the multipole moments are $`m_+,D_+,Q_+`$ etc.. In $`M^{}`$ the local coordinates are $`x_{}^\mu =(x,y,r,u)`$ and the multipole moments are $`m,D,Q`$ etc.. The metric tensors induced on $`u=0`$ (the boundary between $`M^+`$ and $`M^{}`$) by its embedding in $`M^+`$ and $`M^{}`$ agree provided the following matching conditions are satisfied:
$`x_+`$ $`=`$ $`x+2{\displaystyle \frac{[Q]xp_0^1}{r^3}}+\mathrm{},`$ (3.10)
$`y_+`$ $`=`$ $`y+2{\displaystyle \frac{[Q]yp_0^1}{r^3}}+\mathrm{},`$ (3.11)
$`r_+`$ $`=`$ $`r+{\displaystyle \frac{[Q]p_0^2}{r^2}}(x^2+y^22)+\mathrm{}.`$ (3.12)
Applying the Barrabès–Israel technique yields the results given in which in the coordinates $`(x,y,r)`$ read: the components of the stress–energy tensor (2.44) on $`u=0`$ are given by
$`16\pi S^{11}`$ $`=`$ $`16\pi S^{22}={\displaystyle \frac{12[Q]}{r^6}}(x^2+y^22)+\mathrm{},`$ (3.13)
$`16\pi S^{12}`$ $`=`$ $`{\displaystyle \frac{24[Q]}{r^9}}p_0^2xy(x^2+y^22)+\mathrm{},`$ (3.14)
$`16\pi S^{13}`$ $`=`$ $`{\displaystyle \frac{6[D]x}{r^4}}+{\displaystyle \frac{6[Q]}{r^5}}xp_0^1(x^2+y^25)+\mathrm{},`$ (3.15)
$`16\pi S^{23}`$ $`=`$ $`{\displaystyle \frac{6[D]y}{r^4}}+{\displaystyle \frac{6[Q]}{r^5}}yp_0^1(x^2+y^25)+\mathrm{},`$ (3.16)
and
$`16\pi S^{33}={\displaystyle \frac{4[m]}{r^2}}+{\displaystyle \frac{3[D]}{r^3}}p_0^1(x^2+y^24)`$ (3.17)
$`{\displaystyle \frac{3[Q]}{2r^4}}p_0^2\{24+(x^2+y^2)(x^2+y^220)\}+\mathrm{}.`$
The coefficients of the delta function in the Weyl tensor are
$$\widehat{\mathrm{\Psi }}_2=\frac{2[Q]}{r^4}p_0^2(x^2+y^22)+\mathrm{},$$
(3.18)
$$\widehat{\mathrm{\Psi }}_3=\frac{3[D]}{\sqrt{2}r^3}p_0^1(xiy)\frac{3[Q]}{\sqrt{2}r^4}p_0^2(xiy)(x^2+y^25)+\mathrm{},$$
(3.19)
and
$$\widehat{\mathrm{\Psi }}_4=\frac{9[Q]}{4r^4}p_0^2(xiy)^2+\mathrm{}.$$
(3.20)
In (3.18)–(3.20) we have an unconventional peeling behavior. It is clear now by comparing the conditions (3.6)– (3.9) and the general formulas (2.56)–(2.59) that this case is unconventional because the space–times $`M^+`$ and $`M^{}`$ tend to flatness faster, as future null infinity is approached, than the more general space-times $`M^+`$ and $`M^{}`$ considered in section 2. We note that this signal is free of directional singularities. To see unambiguously why $`\widehat{\mathrm{\Psi }}_4`$ describes the impulsive gravitational wave part of the signal and $`\widehat{\mathrm{\Psi }}_2`$ and $`\widehat{\mathrm{\Psi }}_3`$ describe the light–like shell of matter (neutrino burst, for example) the reader must consult .
## 4 Discussion
We noted that (2.56)–(2.59) and (3.3) and (3.4) exhibit the conventional peeling behavior. Normally one associates the radiative part of the field ($`\mathrm{\Psi }_4`$ throughout this paper) with a dominant $`\frac{1}{r}`$–behavior. We see in (3.20) that this is not the case in the multipole example. On the other hand the radiative part (3.20) of that signal is due primarily to the jump in the quadrupole moment of the source across the light–like signal and this is something that would be expected. The general formulas (2.56)–(2.59) allow plenty of scope to construct further examples with unconventional peeling behavior. For example we could take $`M^+`$ and $`M^{}`$ to be both Schwarzschild space– times with masses $`m_+`$ and $`m`$. These space–times can be attached on $`u=0`$ asymptotically (for large $`r`$) with the matching (2.27)–(2.30) with $`f=x,g=y,h=1`$ and with (2.34) and (2.35) holding but with $`\alpha _1=\beta _1=0`$. Now the stress–energy tensor on $`u=0`$ is given by (2.47)–(2.55) with $`a_2=b_2=\alpha _1=\beta _1=0`$. Thus for example
$$S^{33}=\frac{1}{4\pi r^2}\left([m]\frac{1}{2}h_0\right)+\mathrm{},$$
(4.1)
in this case. With $`\widehat{\mathrm{\Psi }}_A`$ calculated from (2.56)– (2.59) we see that in particular
$$\widehat{\mathrm{\Psi }}_4=\frac{2}{r^2}\frac{}{z}\left(p_0^2\frac{h_0}{z}\right)+O\left(\frac{1}{r^3}\right).$$
(4.2)
In general this wave will exhibit directional singularities. For example if $`f_1`$ and $`g_1`$ are constants then
$$\widehat{\mathrm{\Psi }}_4=\frac{1}{8r^2}(f_1ig_1)\overline{z}+O\left(\frac{1}{r^3}\right),$$
(4.3)
with $`z=x+iy`$.
AcknowledgementThis paper was motivated by a stimulating correspondence with Professor J. N. Goldberg. One of us (G. F. B) wishes to thank the Department of Education and Science for a postdoctoral fellowship. |
warning/0002/math0002095.html | ar5iv | text | # Virtual Gromov-Witten Invariants and the Quantum Cohomology Rings of General Type Projective Hypersurfaces
## 1 Introduction and Statement of the Main Results
Recently, some works on the quantum cohomology ring of the general type projective hypersurface have appeared , , . In , Lian, Liu and Yau generalized their mirror principle to the case of the general type projective hypersurface and proposed a theoretical recipe to construct the generating function of a certain type of Gromov-Witten invariants including gravitational descendants from the hypergeometric data. In , Gathmann considered the relative Gromov-Witten invariants of the projective space with various tangency conditions on the ample hypersurface in the projective space. He verified the recursive formulas which increase the tangency condition like the results of Caporaso and Harris , and proposed an algorithm of computing the Gromov-Witten invariants of the hypersurface as a limit of iterative application of the recursive formulas.
In this paper, we continue the analysis of on the quantum Kähler sub-ring $`QH_e^{}(M_N^k)`$, where $`M_N^k`$ is the degree $`k`$ hypersurface in $`CP^{N1}`$, especially in the case when its first Chern class is negative. Our approach is different from the ones of and , mainly because we don’t use Mumford-Morita class (gravitational descendants). In , we proposed the generalized mirror transformation on the quantum cohomology rings of $`M_N^k,(Nk1)`$ that represents the sturctural constant $`L_n^{N,k,d}`$ of $`QH_e^{}(M_N^k)`$ in terms of the virtual sturctural constants $`\stackrel{~}{L}_n^{N,k,d}`$, which is the analogue of the hypergeometric series used in the mirror calculation in our context.
Now we restate the main conjecture in . Let $`P_m`$ be the set of partitions of $`m`$ into positive integers and $`\sigma _m`$ be an element of $`P_m`$. We also denote the length of a partition $`\sigma _m`$ by $`l(\sigma _m)`$ (i.e., $`\sigma _m:m=d_1+d_2+\mathrm{}+d_{l(\sigma _m)},d_1d_2\mathrm{}d_{l(\sigma _m)}1`$). We denote by $`\text{mul}(i,\sigma _m),(1im)`$ the multiplicity of $`i`$ in $`\sigma _m`$. Our previous conjecture is the following:
###### Conjecture 1
The generalized mirror transformation takes the form
$`L_n^{N,k,d}`$ $`=`$ $`{\displaystyle \underset{m=0}{\overset{d1}{}}}{\displaystyle \underset{\sigma _mP_m}{}}(1)^{l(\sigma _m)}{\displaystyle \frac{d^{l(\sigma _m)}}{_{j=1}^{l(\sigma _m)}d_j_{i=1}^m\text{mul}(i,\sigma _m)!}}{\displaystyle \underset{i=1}{\overset{l(\sigma _m)}{}}}\stackrel{~}{L}_{1+(kN)d_i}^{N,k,d_i}G_{dm}^{N,k,d}(n;\sigma _m),`$
(1.1)
where $`G_{dm}^{N,k,d}(n;\sigma _m)`$ is a polynomial of $`\stackrel{~}{L}_n^{N,k,d}`$ with weighted degree $`d`$.
Now, we propose a conjecture on the explicit form of $`G_{dm}^{N,k,d}(n;\sigma _m)`$.
###### Conjecture 2
$`V_{dm}^{N,k,d}(n;\sigma _m)`$ $`:=`$ $`{\displaystyle \frac{1}{k}}v(𝒪_{e^{N2n}}𝒪_{e^{n1(kN)d}}{\displaystyle \underset{i=1}{\overset{l(\sigma _m)}{}}}𝒪_{e^{1+(kN)d_i}})_{dm}{\displaystyle \frac{1}{(dm)^{l(\sigma _m)1}}}`$ (1.2)
$``$ $`G_{dm}^{N,k,d}(n;\sigma _m)`$
where $`v(𝒪_{e^{N2n}}𝒪_{e^{n1(kN)d}}_{i=1}^{l(\sigma _m)}𝒪_{e^{1+(kN)d_i}})_{dm}`$ is the virtual Gromov-Witten invariant defined below. If $`ł(\sigma _m)1`$ or $`dm=1`$, (1.2) becomes an exact equality.
We have to make some remarks on the meaning of $``$ in (1.2). In the case when $`l(\sigma _m)1`$ or $`dm=1`$, we can see the coincidence between $`G_{dm}^{N,k,d}(n;\sigma _m)`$ and $`V_{dm}^{N,k,d}(n;\sigma _m)`$, but when $`l(\sigma _m)2`$ and $`dm2`$, we have to modify slightly the coefficients of some polynomials in $`\stackrel{~}{L}_n^{N,k,d}`$ that appear in $`V_{dm}^{N,k,d}(n;\sigma _m)`$. These cases indeed appear when $`d4`$. Even in such situations, $`V_{dm}^{N,k,d}(n;\sigma _m)`$ strongly predicts the form of $`G_{dm}^{N,k,d}(n;\sigma _m)`$.
Now we turn into the definition of the virtual Gromov-Witten invariants, that are the main ingredients of this paper.
###### Definition 1
The virtual Gromov-Witten invariant $`v(_{j=1}^n𝒪_{e^{a_j}})_d`$ on $`M_N^k`$ is the rational number that satisfy the condition:
(i) initial condition
$`v(𝒪_{e^a}𝒪_{e^b}𝒪_{e^c})_0=k\delta _{a+b+c,N2},`$
$`v({\displaystyle \underset{j=1}{\overset{n}{}}}𝒪_{e^{a_j}})_0=0,(n3),`$
$`{\displaystyle \frac{1}{k}}v(𝒪_{e^{N2n}}𝒪_{e^{n1(kN)d}}𝒪_e)_d=\stackrel{~}{L}_n^{N,k,d}\stackrel{~}{L}_{1+(kN)d}^{N,k,d},(d1),`$ (1.3)
(ii) flat metric condition
$`v(𝒪_{e^0}𝒪_{e^a}𝒪_{e^b})_0=k\delta _{a+b,N2},`$
$`v(𝒪_{e^0}{\displaystyle \underset{j=1}{\overset{n}{}}}𝒪_{e^{a_j}})_d=0,(d1,\text{or}d=0,n2),`$ (1.4)
(iii) topological selection rule
$$v(\underset{j=1}{\overset{n}{}}𝒪_{e^{a_j}})_d0(N5)+(Nk)d=\underset{j=1}{\overset{n}{}}(a_j1),$$
(1.5)
(iv) Kähler equation
$$v(𝒪_e\underset{j=1}{\overset{n}{}}𝒪_{e^{a_j}})_d=dv(\underset{j=1}{\overset{n}{}}𝒪_{e^{a_j}})_d,$$
(1.6)
(v) associativity equation ,
$`{\displaystyle \underset{d_1=0}{\overset{d}{}}}{\displaystyle \underset{\{\alpha _{}\}{\scriptscriptstyle \{\beta _{}\}}=\{n_{}\}}{}}{\displaystyle \underset{i=0}{\overset{N2}{}}}v(𝒪_{e^a}𝒪_{e^b}({\displaystyle \underset{\alpha _j\{\alpha _{}\}}{}}𝒪_{e^{\alpha _j}})𝒪_{e^i})_{d_1}v(𝒪_{e^{N2i}}({\displaystyle \underset{\beta _j\{\beta _{}\}}{}}𝒪_{e^{\beta _j}})𝒪_{e^c}𝒪_{e^d})_{dd_1}`$
$`={\displaystyle \underset{d_1=0}{\overset{d}{}}}{\displaystyle \underset{\{\alpha _{}\}{\scriptscriptstyle \{\beta _{}\}}=\{n_{}\}}{}}{\displaystyle \underset{i=0}{\overset{N2}{}}}v(𝒪_{e^a}𝒪_{e^c}({\displaystyle \underset{\alpha _j\{\alpha _{}\}}{}}𝒪_{e^{\alpha _j}})𝒪_{e^i})_{d_1}v(𝒪_{e^{N2i}}({\displaystyle \underset{\beta _j\{\beta _{}\}}{}}𝒪_{e^{\beta _j}})𝒪_{e^b}𝒪_{e^d})_{dd_1},`$
$`(a+b+c+d+{\displaystyle \underset{j=1}{\overset{m}{}}}(n_j1)=N2+(Nk)d).`$ (1.7)
The difference between the virtual Gromov-Witten invariants and the ordinary Gromov-Witten invariants comes from the initial condition (1.3). In ordinary cases, we set
$$\frac{1}{k}𝒪_{e^{N2n}}𝒪_{e^{n1(kN)d}}𝒪_e_d=L_n^{N,k,d},(d1).$$
(1.8)
(1.3) and (1.8) indeed differ when $`d>1`$. With the above conditions, we can completely determine the virtual Gromov-Witten invariants like the ordinary Gromov-Witten invariants of genus $`0`$. Using the generating function of the virtual Gromov-Witten invariants:
$`F_v(z,t_i)={\displaystyle \frac{1}{3!}}{\displaystyle \underset{i,j,k=0}{\overset{N2}{}}}k\delta _{i+j+k,N2}t_it_jt_k+f_v(z,t_i),`$ (1.9)
we can state the main conjecture in a more compact form,
$$L_n^{N,k,d}\frac{1}{k}_{C_0}𝑑z\frac{1}{z^{d+1}}\mathrm{exp}(d\underset{j=1}{\overset{\mathrm{}}{}}\frac{\stackrel{~}{L}_{1+(kN)j}^{N,k,j}}{j}\frac{_{1+(kN)j}}{_1}z^j)_{N2n}_{n1d(kN)}_1f_v(z,t_i)|_{t_i=0},$$
(1.10)
where $`t_j(j=0,\mathrm{}N2)`$ is the variable that couples to the element $`𝒪_{e^j}`$ of $`QH_e^{}(M_N^k)`$ and $`z`$ is the formal degree counting variable.
This paper is organized as follows. In Section 2, we introduce the notation of the quantum Kähler subring of $`M_N^k`$, and review the results obtained in , and . In Section 3, we reproduce the formulas obtained in under the assumption of Conjecture 2. In Section 4, we construct the explicit form of the generalized mirror transformation for degree $`4,5`$ rational Gromov -Witten invariants of $`M_{k1}^k`$ using the Conjecture 2 and some numerical data obtained from the fixed point computation in .
## 2 Quantum Kähler Sub-ring of Projective Hypersurfaces
### 2.1 Notation
In this section, we introduce the quantum Kähler sub-ring of the quantum cohomology ring of a degree $`k`$ hypersurface in $`CP^{N1}`$. Let $`M_N^k`$ be a hypersurface of degree $`k`$ in $`CP^{N1}`$. We denote by $`QH_e^{}(M_N^k)`$ the subring of the quantum cohomology ring $`QH^{}(M_N^k)`$ generated by $`𝒪_e`$ induced from the Kähler form $`e`$ (or, equivalently the intersection $`HM_N^k`$ between a hyperplane class $`H`$ of $`CP^{N1}`$ and $`M_N^k`$). Additive basis of $`QH_e^{}(M_N^k)`$ is given by $`𝒪_{e^j}(j=0,1,\mathrm{},N2)`$, which is induced from $`e^jH^{j,j}(M_N^k)`$. The multiplication rule of $`QH_e^{}(M_N^k)`$ is determined by the Gromov-Witten invariant of genus $`0`$ $`𝒪_e𝒪_{e^{N2m}}𝒪_{e^{m1(kN)d}}_{d,M_N^k}`$ and it is given as follows:
$`L_m^{N,k,d}`$ $`:=`$ $`{\displaystyle \frac{1}{k}}𝒪_e𝒪_{e^{N2m}}𝒪_{e^{m1(kN)d}}_d,`$
$`𝒪_e1`$ $`=`$ $`𝒪_e,`$
$`𝒪_e𝒪_{e^{N2m}}`$ $`=`$ $`𝒪_{e^{N1m}}+{\displaystyle \underset{d=1}{\overset{\mathrm{}}{}}}L_m^{N,k,d}q^d𝒪_{e^{N1m+(kN)d}},`$
$`q`$ $`:=`$ $`\mathrm{exp}(t),`$ (2.11)
where the subscript $`d`$ counts the degree of the rational curves measured by $`e`$. Therefore, $`q=\mathrm{exp}(t)`$ is the degree counting parameter.
###### Definition 2
We call $`L_n^{N,k,d}`$ the structural constant of weighted degree $`d`$.
Since $`M_N^k`$ is a complex $`(N2)`$ dimensional manifold, we see that a structure constant $`L_m^{N,k,d}`$ is non-zero only if the following condition is satisfied:
$`1N2mN2,1m1+(Nk)dN2,`$ (2.12)
$``$ $`max\{0,2(Nk)d\}mmin\{N3,N1(Nk)d\}.`$
We rewrite (2.12) into
$`(Nk2)`$ $``$ $`0m(N1)(Nk)d`$
$`(Nk=1,d=1)`$ $``$ $`1mN3`$
$`(Nk=1,d2)`$ $``$ $`0mN1(Nk)d`$
$`(Nk0)`$ $``$ $`2+(kN)dmN3.`$ (2.13)
From (2.13), we easily see that the number of the non-zero structure constants $`L_m^{N,k,d}`$ is finite except for the case of $`N=k`$. Moreover, if $`N2k`$, the non-zero structure constants come only from the $`d=1`$ part and the non-vanishing $`L_m^{N,k,1}`$ is determined by $`k`$ and independent of $`N`$. The $`N2k`$ region is studied by Beauville , and his result plays the role of an initial condition of our discussion later. Explicitly, they are given by the formula :
$$\underset{n=0}{\overset{k1}{}}L_n^{N,k,1}w^n=k\underset{j=1}{\overset{k1}{}}(jw+(kj)),$$
(2.14)
and the other $`L_n^{N,k,d}`$’s all vanishes. In the case of $`N=k`$, the multiplication rule of $`QH_e^{}(M_k^k)`$ is given as follows:
$`𝒪_e1`$ $`=`$ $`𝒪_e,`$
$`𝒪_e𝒪_{e^{k2m}}`$ $`=`$ $`(1+{\displaystyle \underset{d=1}{\overset{\mathrm{}}{}}}q^dL_m^{k,k,d})𝒪_{e^{k1m}}(m=2,3,\mathrm{},k3),`$
$`𝒪_e𝒪_{e^{k3}}`$ $`=`$ $`𝒪_{e^{k2}}.`$ (2.15)
We introduce here the generating function of the structure constants of the Calabi-Yau hypersurface $`M_k^k`$:
$$L_m^{k,k}(e^t):=1+\underset{d=1}{\overset{\mathrm{}}{}}L_m^{k,k,d}e^{dt}(m=2,\mathrm{},k3).$$
(2.16)
### 2.2 Review of Results for Fano and Calabi-Yau Hypersurfaces and Virtual Structure Constants
Let us summarize the results of , . In , we showed that the structure constants $`L_m^{N,k,d}`$ of $`QH_e^{}(M_N^k)`$ for $`(Nk2)`$ can be obtained by applying the recursive formulas which describe $`L_m^{N,k,d}`$ in terms of $`L_m^{}^{N+1,k,d^{}}(d^{}d)`$, with the initial conditions of $`L_m^{N,k,1}`$ given by (2.14) and $`L_m^{N,k,d}=0(d2)`$ in the $`N2k`$ region. Let us introduce the construction of the recursive formulas given in . First, we introduce the polynomial $`Poly_d`$ in $`x,y,z_1,z_2,\mathrm{},z_{d1}`$ defined by the formula:
$`Poly_d(x,y,z_1,z_2,\mathrm{},z_{d1})`$
$`={\displaystyle \frac{1}{(2\pi \sqrt{1})^{d1}}}{\displaystyle _{C_1}}{\displaystyle \frac{dt_1}{t_1}}\mathrm{}{\displaystyle _{C_{d1}}}{\displaystyle \frac{dt_{d1}}{t_{d1}}}{\displaystyle \underset{j=1}{\overset{d1}{}}}({\displaystyle \frac{(dj)x+jy}{d}}+{\displaystyle \underset{i=1}{\overset{j}{}}}{\displaystyle \frac{dj}{di}}t_i+{\displaystyle \underset{i=j+1}{\overset{d1}{}}}{\displaystyle \frac{j}{i}}t_i+`$
$`z_j({\displaystyle \frac{(dj)x+jy}{d}}+{\displaystyle \underset{i=1}{\overset{j}{}}}{\displaystyle \frac{dj}{di}}t_i+{\displaystyle \underset{i=j+1}{\overset{d1}{}}}{\displaystyle \frac{j}{i}}t_i)/({\displaystyle \frac{(dj)x+jy}{d}}+{\displaystyle \underset{i=1}{\overset{j}{}}}{\displaystyle \frac{dj}{di}}t_i+{\displaystyle \underset{i=j+1}{\overset{d1}{}}}{\displaystyle \frac{j}{i}}t_iz_j)).`$
(2.17)
In (2.17), we have to choose the path $`C_i`$ carefully to obtain the correct answer. See for details. Consider the monomial $`x^{d_{i_0}}z_{i_1}^{d_{i_1}}\mathrm{}z_{i_m}^{d_{i_m}}y^{d_{i_{m+1}}}(_{j=0}^{m+1}d_{i_j}=d1)`$, that appear in $`Poly_d`$, associated with the following ordered partition of a positive integer $`d`$ :
$$0=i_0<i_1<i_2<\mathrm{}<i_m<i_{m+1}=d.$$
(2.18)
Next, we prepare some elements in (a free abelian group) $`𝐙^{m+1}`$, which are determined for each monomial $`x^{d_{i_0}}z_{i_1}^{d_{i_1}}\mathrm{}z_{i_m}^{d_{i_m}}y^{d_{i_{m+1}}}`$ , as follows:
$`\alpha `$ $`:=`$ $`(m+1d,m+1d,\mathrm{},m+1d),`$
$`\beta `$ $`:=`$ $`(0,i_11,i_22,\mathrm{},i_mm),`$
$`\gamma `$ $`:=`$ $`(0,i_1(Nk),i_2(Nk),\mathrm{},i_m(Nk)),`$
$`ϵ_1`$ $`:=`$ $`(1,0,0,0,\mathrm{},0),`$
$`ϵ_2`$ $`:=`$ $`(1,1,0,0,\mathrm{},0),`$
$`ϵ_3`$ $`:=`$ $`(1,1,1,0,\mathrm{},0),`$
$`\mathrm{}`$
$`ϵ_{m+1}`$ $`:=`$ $`(1,1,1,1,\mathrm{},1).`$ (2.19)
Now we define $`\delta =(\delta _1,\mathrm{},\delta _{m+1})𝐙^{m+1}`$ by the formula:
$$\delta :=\alpha +\beta +\gamma +\underset{j=1}{\overset{m}{}}(d_{i_j}1)ϵ_j+d_{i_{m+1}}ϵ_{m+1}.$$
(2.20)
Then the recursive formulas are given as follows:
$$L_n^{N,k,d}=\varphi (Poly_d),$$
(2.21)
where $`\varphi `$ is a $`𝐐`$-linear map from the $`𝐐`$-vector space of the homogeneous polynomials of degree $`d1`$ in $`x,y,z_1,\mathrm{},z_{d1}`$ to the $`𝐐`$-vector space of the weighted homogeneous polynomials of degree $`d`$ in $`L_m^{N+1,k,d^{}}`$. And it is defined on the basis by:
$$\varphi (x^{d_0}y^{d_d}z_{i_1}^{d_{i_1}}\mathrm{}z_{i_m}^{d_{i_m}})=\underset{j=1}{\overset{m+1}{}}L_{n+\delta _j}^{N+1,k,i_ji_{j1}}.$$
(2.22)
In the $`d5`$ cases, we examined that these recursive formulas naturally lead us to the relation:
$$(𝒪_e)^{N1}k^k(𝒪_e)^{k1}q=0,$$
(2.23)
of $`QH_e^{}(M_N^k)(Nk2)`$ by descending induction using Beauville’s result , , . In the $`Nk=1`$ case, the recursive formulas receive modification only in the $`d=1`$ part:
$$L_m^{k+1,k,1}=L_m^{k+2,k,1}L_0^{k+2,k,1}=L_m^{k+2,k,1}k!.$$
(2.24)
This leads us to the following relation of $`QH_e^{}(M_{k+1}^k)`$:
$$(𝒪_e+k!q)^{N1}k^k(𝒪_e+k!q)^{k1}q=0.$$
(2.25)
The structural constant $`L_m^{k,k,d}`$ for a Calabi-Yau hypersurface does not obey the recursive formulas. Instead, we introduce here the virtual structure constants $`\stackrel{~}{L}_m^{N,k,d}`$ as follows.
###### Definition 3
Let $`\stackrel{~}{L}_m^{N,k,d}`$ be the rational number obtained by applying the recursion relations of Fano hypersurfaces for arbitrary $`N`$ and $`k`$ with the initial condition $`L_n^{N,k,1}(N2k)`$ and $`L_n^{N,k,d}=0(d2,N2k)`$.
###### Remark 1
In the $`Nk2`$ region, $`\stackrel{~}{L}_m^{N,k,d}=L_m^{N,k,d}`$.
###### Definition 4
We call $`\stackrel{~}{L}_n^{N,k,d}`$ the virtual structural constant of weighted degree $`d`$.
We define the generating function of the virtual structural constants of the Calabi-Yau hypersurface $`M_k^k`$ as follows:
$`\stackrel{~}{L}_n^{k,k}(e^x)`$ $`:=`$ $`1+{\displaystyle \underset{d=1}{\overset{\mathrm{}}{}}}\stackrel{~}{L}_n^{k,k,d}e^{dx},`$ (2.26)
$`(n=0,1,\mathrm{},k1).`$
In , we observed that $`\stackrel{~}{L}_n^{k,k}(e^x)`$ gives us the information of the B-model of the mirror manifold of $`M_k^k`$. More explicitly, we conjectured
$`\stackrel{~}{L}_0^{k,k}(e^x)={\displaystyle \underset{d=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(kd)!}{(d!)^k}}e^{dx},`$
$`\stackrel{~}{L}_1^{k,k}(e^x)={\displaystyle \frac{dt(x)}{dx}}:={\displaystyle \frac{d}{dx}}\left(x+({\displaystyle \underset{d=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(kd)!}{(d!)^k}}({\displaystyle \underset{i=1}{\overset{d}{}}}{\displaystyle \underset{m=1}{\overset{k1}{}}}{\displaystyle \frac{m}{i(kim)}})e^{dx})/({\displaystyle \underset{d=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(kd)!}{(d!)^k}}e^{dx})\right)`$ (2.27)
where the r.h.s. of (2.27) is derived from the solutions of the ODE for the period integral of the mirror manifold of $`M_k^k`$,
$$((\frac{d}{dx})^{k1}ke^x(k\frac{d}{dx}+1)(k\frac{d}{dx}+2)\mathrm{}(k\frac{d}{dx}+k1))w(x)=0,$$
(2.28)
that was used in the computation based on the mirror symmetry, . Of course, we can extend the conjecture (2.27) to the general $`\stackrel{~}{L}_n^{k,k}(e^x)`$ if we compare the $`\stackrel{~}{L}_n^{k,k}(e^x)`$ with the B-model three point functions in . Hence we obtain the mirror map $`t=t(x)`$ without using the mirror conjecture:
$$t(x)=x+_{\mathrm{}}^x𝑑x^{}(\stackrel{~}{L}_1^{k,k}(e^x^{})1)=x+\underset{d=1}{\overset{\mathrm{}}{}}\frac{\stackrel{~}{L}_1^{k,k,d}}{d}e^{dx}.$$
(2.29)
With the conjecture given by (2.27), we can construct the mirror transformation that transforms the virtual structural constants of the Calabi-Yau hypersurface into the real ones as follows:
$$L_m^{k,k}(e^t)=\frac{\stackrel{~}{L}_m^{k,k}(e^{x(t)})}{\stackrel{~}{L}_1^{k,k}(e^{x(t)})}.(m=2,\mathrm{},k3)$$
(2.30)
The above formula is further rewritten as follows:
$$L_n^{k,k,d}=\underset{m=0}{\overset{d1}{}}\text{Res}_{z=0}(z^{m1}\mathrm{exp}(d\underset{j=1}{\overset{\mathrm{}}{}}\frac{\stackrel{~}{L}_1^{k,k,j}}{j}z^j))(\stackrel{~}{L}_n^{k,k,dm}\stackrel{~}{L}_1^{k,k,dm}).$$
(2.31)
This formula motivated us to propose the conjecture described in (1.1) or in (1.10) .
## 3 Derivation of the Previous Results
In this section, we first show that $`V_{dm}^{N,k,d}(n;\sigma _m)`$ satisfy the ansatz of $`G_{dm}^{N,k,d}(n;\sigma _m)`$ proposed in . Then, we show that $`V_{dm}^{N,k,d}(n;\sigma _m)`$ coincides with $`G_{dm}^{N,k,d}(n;\sigma _m)`$ in the $`d3`$ cases.
###### Proposition 1
(i) flat metric condition
$$V_{dm}^{N,k,d}(1+(kN)d;\sigma _m)=V_{dm}^{N,k,d}(N2;\sigma _m)=0.$$
(3.32)
(ii) symmetry
$$V_{dm}^{N,k,d}(n;\sigma _m)=V_{dm}^{N,k,d}(N1+(kN)dn;\sigma _m).$$
(3.33)
(iii)
$$V_d^{N,k,d}(n;(0))=\stackrel{~}{L}_n^{N,k,d}\stackrel{~}{L}_{1+d(kN)}^{N,k,d}.$$
(3.34)
(iv)
$$V_{dm}^{N,k,d}(2+(kN)(d+f);\sigma _m)=V_{dm}^{N,k,d+f}(2+(kN)(d+f);\sigma _m(f)).$$
(3.35)
proof)
(i), (ii) and (iii) are obvious by definition. Here, we give a proof of (iv). By the definition of $`V_{dm}^{N,k,d+f}(n;\sigma _m(f))`$ and by Kähler equation, we have
$`V_{dm}^{N,k,d+f}(2+(kN)(d+f);\sigma _m(f))`$
$`={\displaystyle \frac{1}{k}}{\displaystyle \frac{1}{(dm)^{l(\sigma _m)}}}v(𝒪_{e^{N22(kN)(d+f)}}𝒪_e𝒪_{e^{1+(kN)f}}{\displaystyle \underset{j=1}{\overset{l(\sigma _m)}{}}}𝒪_{e^{1+(kN)d_j}})_{dm}`$
$`={\displaystyle \frac{1}{k}}{\displaystyle \frac{1}{(dm)^{l(\sigma _m)1}}}v(𝒪_{e^{N22(kN)(d+f)}}𝒪_{e^{2+(kN)(d+f)1(kN)d}}{\displaystyle \underset{j=1}{\overset{l(\sigma _m)}{}}}𝒪_{e^{1+(kN)d_j}})_{dm}.`$
(3.36)
The last line of (3.36) is nothing but $`V_{dm}^{N,k,d}(2+(kN)(d+f);\sigma _m)`$. Q.E.D.
Then we introduce the definition:
###### Definition 5
$$\stackrel{~}{V}_{dm}^{N,k,d+f}(n;\sigma _m(f)):=\underset{j=0}{\overset{(kN)f}{}}V_{dm}^{N,k,d}(nj;\sigma _m)\underset{j=0}{\overset{(kN)f}{}}V_{dm}^{N,k,d}(1+(kN)(d+f)j;\sigma _m).$$
(3.37)
We denote by $`\pi _f`$ the map which maps the function $`g(n)`$ on $`𝐙`$ to
$`_{j=0}^{(kN)f}g(nj)_{j=0}^{(kN)f}g(1+(kN)(d+f)j)`$. Then we have
$$\pi _f(V_{dm}^{N,k,d}(n;\sigma _m))=\stackrel{~}{V}_{dm}^{N,k,d+f}(n;\sigma _m(f)).$$
(3.38)
With this definition, we are led to the following proposition.
###### Proposition 2
$`\stackrel{~}{V}_{dm}^{N,k,d+f}(n;\sigma _m(f))`$ satisfies the condition (iv) of the above proposition:
$$\stackrel{~}{V}_{dm}^{N,k,d+f}(2+(kN)(d+f);\sigma _m(f))=V_{dm}^{N,k,d}(2+(kN)(d+f);\sigma _m),$$
(3.39)
and the condition (i) and (ii).
proof) The fact that $`\stackrel{~}{V}_{dm}^{N,k,d+f}(n;\sigma _m(f))`$ satisfies the condition (i) is rather obvious, and we first prove the condition (ii). It suffices to consider the part $`_{j=0}^{(kN)f}V_{dm}^{N,k,d}(nj;\sigma _m)`$:
$`{\displaystyle \underset{j=0}{\overset{(kN)f}{}}}V_{dm}^{N,k,d}(N1+(kN)(d+f)nj;\sigma _m)`$ (3.40)
$`=`$ $`{\displaystyle \underset{j=0}{\overset{(kN)f}{}}}V_{dm}^{N,k,d}(N1+(kN)dn+j;\sigma _m)`$
$`=`$ $`{\displaystyle \underset{j=0}{\overset{(kN)f}{}}}V_{dm}^{N,k,d}(nj;\sigma _m).`$
Next, we turn to the formula $`(\text{3.39})`$. By definition and by the condition $`V_{dm}^{N,k,d}(1+(kN)d;\sigma _m)=0`$, we obtain
$`\stackrel{~}{V}_{dm}^{N,k,d+f}(2+(kN)(d+f);\sigma _m(f))`$ (3.41)
$`=`$ $`{\displaystyle \underset{j=0}{\overset{(kN)f}{}}}V_{dm}^{N,k,d}(2+(kN)(d+f)j;\sigma _m){\displaystyle \underset{j=0}{\overset{(kN)f}{}}}V_{dm}^{N,k,d}(1+(kN)(d+f)j;\sigma _m)`$
$`=`$ $`V_{dm}^{N,k,d}(2+(kN)(d+f);\sigma _m).`$
Q.E.D.
###### Definition 6
$$hi_{dm}^{N,k,d+f}(n;\sigma _m(f)):=V_{dm}^{N,k,d+f}(n;\sigma _m(f))\stackrel{~}{V}_{dm}^{N,k,d+f}(n;\sigma _m(f))$$
(3.42)
###### Proposition 3
$`hi_{dm}^{N,k,d+f}(n;\sigma _m(f))`$ satisfies the condition:
$$hi_{dm}^{N,k,d+f}(2+(kN)(d+f);\sigma _m(f))=hi_{dm}^{N,k,d+f}(1+(kN)(d+f);\sigma _m(f))=0,$$
(3.43)
and the condition (i), (ii).
proof) Immediate. Q.E.D.
With the above discussions, we can construct the following decomposition of $`V_{dm}^{N,k,d}(n;\sigma _m)(\sigma _m:m=d_1+d_2+\mathrm{}+d_{l(\sigma _m)},d_1\mathrm{}d_{l(\sigma _m)}1)`$:
$`V_{dm}^{N,k,d}(n;\sigma _m)`$ $`=`$ $`\pi _{d_{l(\sigma _m)}}\mathrm{}\pi _{d_2}\pi _{d_1}(\stackrel{~}{L}_n^{N,k,d}\stackrel{~}{L}_{1+(kN)d}^{N,k,d})`$ (3.44)
$`+{\displaystyle \underset{j=2}{\overset{l(\sigma _m)}{}}}\pi _{d_{l(\sigma _m)}}\mathrm{}\pi _{d_j}(hi^{(j)})+hi^{(l(\sigma _m))}.`$
###### Remark 2
The above decomposition depends on the order of $`d_i`$’s. Hence it is not unique.
In (3.44), we omit the subscripts of $`hi_{dm}^{N,k,d+f}(n;\sigma _m(f))`$. We can easily see that $`hi_{dm}^{N,k,d+f}(n;\sigma _m(f))`$ consists of monomials of degree $`(dm)`$ of $`\stackrel{~}{L}_j^{N,k,m^{}}(m^{}<dm)`$ only, because the linear dependence on $`\stackrel{~}{L}_j^{N,k,dm}`$ cannot satisfy the condition (3.43). Thus we are led to the proposition by picking up the top term of the decomposition in (3.44):
###### Proposition 4
The linear part of $`V_{dm}^{N,k,d}(n;\sigma _m)`$ is given by the formula:
$`{\displaystyle \underset{j_1=0}{\overset{d_1(kN)}{}}}{\displaystyle \underset{j_2=0}{\overset{d_2(kN)}{}}}\mathrm{}{\displaystyle \underset{j_{l(\sigma _m)}=0}{\overset{d_{l(\sigma _m)}(kN)}{}}}(\stackrel{~}{L}_{n_{i=1}^{l(\sigma _m)}j_i}^{N,k,dm}\stackrel{~}{L}_{1+(kN)d_{i=1}^{l(\sigma _m)}j_i}^{N,k,dm}).`$
(3.45)
Especially in the $`dm=1`$ case, we obtain the following equality:
$`V_1^{N,k,d}(n;\sigma _{d1})`$ $`=`$ $`{\displaystyle \underset{j_1=0}{\overset{d_1(kN)}{}}}{\displaystyle \underset{j_2=0}{\overset{d_2(kN)}{}}}\mathrm{}{\displaystyle \underset{j_{l(\sigma _{d1})}=0}{\overset{d_{l(\sigma _{d1})}(kN)}{}}}(\stackrel{~}{L}_{n_{i=1}^{l(\sigma _{d1})}j_i}^{N,k,1}\stackrel{~}{L}_{1+(kN)d_{i=1}^{l(\sigma _{d1})}j_i}^{N,k,1}).`$
###### Remark 3
By introducing the polynomial in $`x`$:
$`{\displaystyle \underset{j=1}{\overset{l(\sigma _m)}{}}}{\displaystyle \frac{(1x^{d_j(kN)+1})}{(1x)}}={\displaystyle \underset{j=0}{\overset{(kN)m}{}}}A_j^{N,k}(\sigma _m)x^j,`$ (3.47)
we can write the linear part (3.45) in a more compact form,
$`{\displaystyle \underset{j_1=0}{\overset{d_1(kN)}{}}}{\displaystyle \underset{j_2=0}{\overset{d_2(kN)}{}}}\mathrm{}{\displaystyle \underset{j_{l(\sigma _m)}=0}{\overset{d_{l(\sigma _m)}(kN)}{}}}(\stackrel{~}{L}_{n_{i=1}^{l(\sigma _m)}j_i}^{N,k,dm}\stackrel{~}{L}_{1+(kN)d_{i=1}^{l(\sigma _m)}j_i}^{N,k,dm})`$
$`={\displaystyle \underset{j=0}{\overset{(kN)m}{}}}A_j^{N,k}(\sigma _m)(\stackrel{~}{L}_{nj}^{N,k,dm}\stackrel{~}{L}_{1+(kN)dj}^{N,k,dm}).`$ (3.48)
These are the results directly applying the constraints obtained in Proposition 1. On the other hand, we can explicitly express $`V_{dm}^{N,k,d}(n;\sigma _m)`$ in terms of $`\stackrel{~}{L}_n^{N,k,d}`$, because the virtual Gromov-Witten invariants satisfy the Kähler equation and the associativity equation. As examples, we compute $`V_{dm}^{N,k,d}(n;\sigma _m)(d3)`$.
###### Proposition 5
$`V_{dm}^{N,k,d}(n;\sigma _m)(d3)`$ can be written in terms of $`\stackrel{~}{L}_n^{N,k,d}(d3)`$ as follows.
$`V_1^{N,k,1}(n;(0))`$ $`=`$ $`\stackrel{~}{L}_n^{N,k,1}\stackrel{~}{L}_{1+(kN)}^{N,k,1},`$
$`V_2^{N,k,2}(n;(0))`$ $`=`$ $`\stackrel{~}{L}_n^{N,k,2}\stackrel{~}{L}_{1+2(kN)}^{N,k,2},`$
$`V_1^{N,k,2}(n;(1))`$ $`=`$ $`{\displaystyle \underset{j=0}{\overset{kN}{}}}(\stackrel{~}{L}_{nj}^{N,k,1}\stackrel{~}{L}_{1+2(kN)j}^{N,k,1}),`$
$`V_3^{N,k,3}(n;(0))`$ $`=`$ $`\stackrel{~}{L}_n^{N,k,3}\stackrel{~}{L}_{1+3(kN)}^{N,k,3},`$
$`V_2^{N,k,3}(n;(1))`$ $`=`$ $`{\displaystyle \underset{j=0}{\overset{kN}{}}}(\stackrel{~}{L}_{nj}^{N,k,2}\stackrel{~}{L}_{1+3(kN)j}^{N,k,2})+hi_2^{N,k,3}(n;(1))`$
$`V_1^{N,k,3}(n;(2))`$ $`=`$ $`{\displaystyle \underset{j=0}{\overset{2(kN)}{}}}(\stackrel{~}{L}_{nj}^{N,k,1}\stackrel{~}{L}_{1+3(kN)j}^{N,k,1}),`$
$`V_1^{N,k,3}(n;(1)+(1))`$ $`=`$ $`{\displaystyle \underset{j=0}{\overset{2(kN)}{}}}A_j^{N,k}((1)+(1))(\stackrel{~}{L}_{nj}^{N,k,1}\stackrel{~}{L}_{1+3(kN)j}^{N,k,1}).`$ (3.49)
where $`hi_2^{N,k,3}(n;(1))`$ is given by the formula:
$`hi_2^{N,k,3}(n;(1))`$ $`=`$ $`{\displaystyle \underset{j=0}{\overset{(kN)1}{}}}({\displaystyle \underset{m=0}{\overset{j}{}}}\stackrel{~}{L}_{nm}^{N,k,1}\stackrel{~}{L}_{n2(kN)+jm}^{N,k,1}\stackrel{~}{L}_{(kN)+2+j}^{N,k,1}{\displaystyle \underset{m=0}{\overset{2(kN)}{}}}\stackrel{~}{L}_{nm}^{N,k,1}`$ (3.50)
$`+\stackrel{~}{L}_{1+(kN)}^{N,k,1}{\displaystyle \underset{m=j+1}{\overset{2(kN)j1}{}}}\stackrel{~}{L}_{nm}^{N,k,1})`$
$`{\displaystyle \underset{j=0}{\overset{(kN)1}{}}}({\displaystyle \underset{m=0}{\overset{j}{}}}\stackrel{~}{L}_{1+3(kN)m}^{N,k,1}\stackrel{~}{L}_{1+(kN)+jm}^{N,k,1}\stackrel{~}{L}_{(kN)+2+j}^{N,k,1}{\displaystyle \underset{m=0}{\overset{2(kN)}{}}}\stackrel{~}{L}_{1+3(kN)m}^{N,k,1}`$
$`+\stackrel{~}{L}_{1+(kN)}^{N,k,1}{\displaystyle \underset{m=j+1}{\overset{2(kN)j1}{}}}\stackrel{~}{L}_{1+3(kN)m}^{N,k,1}).`$
proof)
We give a proof of the formula for $`V_2^{N,k,3}(n;(1))`$, because the other formulas follow obviously from the preceding discussions. First, we introduce the following virtual G-W invariant:
$`v(𝒪_{e^{N2n}}𝒪_{e^{n12(kN)m}}𝒪_{e^{1+m}})_2.`$ (3.51)
Using the associativity equation, we obtain the equality,
$`{\displaystyle \frac{1}{k}}v(𝒪_{e^{N2n}}𝒪_{e^{n12(kN)m}}𝒪_{e^{1+m}})_2{\displaystyle \frac{1}{k}}v(𝒪_{e^{N2n}}𝒪_{e^{n2(kN)m}}𝒪_{e^m})_2`$
$`=\stackrel{~}{L}_{nm}^{N,k,2}\stackrel{~}{L}_{1+m+2(kN)}^{N,k,2}`$
$`+{\displaystyle \underset{j=0}{\overset{m1}{}}}(\stackrel{~}{L}_{nj}^{N,k,1}\stackrel{~}{L}_{1+(kN)+j}^{N,k,1})(\stackrel{~}{L}_{nm(kN)}^{N,k,1}\stackrel{~}{L}_{1+(kN)}^{N,k,1})`$
$`{\displaystyle \underset{j=0}{\overset{m+(kN)}{}}}(\stackrel{~}{L}_{nj}^{N,k,1}\stackrel{~}{L}_{1+(kN)+j}^{N,k,1})(\stackrel{~}{L}_{1+m+(kN)}^{N,k,1}\stackrel{~}{L}_{1+(kN)}^{N,k,1}).`$ (3.52)
Adding up (3.52) with $`m=0,1,\mathrm{},(kN)`$ and using some algebras, we can reach the desired formula in the proposition. Q.E.D.
From this proposition and the conjectural form of the generalized mirror transformation for the $`d3`$ rational G-W invariants in , we can see the equality $`G_{dm}^{N,k,d}(n;\sigma _m)=V_{dm}^{N,k,d}(n;\sigma _m)(d3)`$. Thus, we have derived the generalized mirror transformation up to $`d3`$ cases in under the assumption of Conjecture 2.
## 4 Explicit Determination in the $`kN=1,d=4,5`$ cases
In this section, we restrict $`kN`$ to $`1`$ to avoid the complication of the formulas. In this setting, we compute $`V_{dm}^{k1,k,d}(n;\sigma _m)`$ in the $`d=4,5`$ cases and discuss the modification of $`V_{dm}^{k1,k,d}(n;\sigma _m)`$ into $`G_{dm}^{k1,k,d}(n;\sigma _m)`$. With the aid of some numerical data, we fix the modification and derive the generalized mirror transformation in these cases. Generalization to the general $`kN`$ cases is rather straightforward.
First, we repeatedly use the associativity equation and obtain the following formula that represent $`V_{4m}^{k1,k,4}(n;\sigma _m)`$ in terms of $`\stackrel{~}{L}_n^{k1,k,d}`$.
###### Proposition 6
$`V_{4m}^{k1,k,4}(n;\sigma _m)`$’s are inductively determined as follows:
$`V_4^{k1,k,4}(n;(0))`$ $`=`$ $`\stackrel{~}{L}_n^{k1,k,4}\stackrel{~}{L}_5^{k1,k,4},`$
$`V_3^{k1,k,4}(n;(1))`$ $`=`$ $`\stackrel{~}{L}_n^{k1,k,3}+\stackrel{~}{L}_{n1}^{k1,k,3}\stackrel{~}{L}_5^{k1,k,3}\stackrel{~}{L}_4^{k1,k,3}`$
$`+(\stackrel{~}{L}_n^{k1,k,2}\stackrel{~}{L}_3^{k1,k,2})(\stackrel{~}{L}_{n3}^{k1,k,1}\stackrel{~}{L}_2^{k1,k,1})`$
$`+(\stackrel{~}{L}_n^{k1,k,1}\stackrel{~}{L}_2^{k1,k,1})(\stackrel{~}{L}_{n2}^{k1,k,2}\stackrel{~}{L}_3^{k1,k,2})`$
$`(\stackrel{~}{L}_3^{k1,k,1}\stackrel{~}{L}_2^{k1,k,1})V_2^{k1,k,4}(n;(2))`$
$`V_1^{k1,k,4}(n;(3))(\stackrel{~}{L}_4^{k1,k,2}\stackrel{~}{L}_3^{k1,k,2}),`$
$`V_2^{k1,k,4}(n;(2))`$ $`=`$ $`V_2^{k1,k,3}(n;(1))+\stackrel{~}{L}_{n2}^{k1,k,2}\stackrel{~}{L}_5^{k1,k,2}`$
$`+V_1^{k1,k,3}(n;(2))(\stackrel{~}{L}_{n3}^{k1,k,1}\stackrel{~}{L}_2^{k1,k,1})`$
$`V_1^{k1,k,4}(n;(3))(\stackrel{~}{L}_4^{k1,k,1}\stackrel{~}{L}_2^{k1,k,1}),`$
$`V_1^{k1,k,4}(n;(3))`$ $`=`$ $`\stackrel{~}{L}_n^{k1,k,1}+\stackrel{~}{L}_{n1}^{k1,k,1}+\stackrel{~}{L}_{n2}^{k1,k,1}+\stackrel{~}{L}_{n3}^{k1,k,1}`$
$`(\stackrel{~}{L}_5^{k1,k,1}+\stackrel{~}{L}_4^{k1,k,2}+\stackrel{~}{L}_3^{k1,k,2}+\stackrel{~}{L}_2^{k1,k,1}),`$
$`V_1^{k1,k,4}(n;(1)+(2))`$ $`=`$ $`\stackrel{~}{L}_n^{k1,k,1}+2\stackrel{~}{L}_{n1}^{k1,k,1}+2\stackrel{~}{L}_{n2}^{k1,k,1}+\stackrel{~}{L}_{n3}^{k1,k,1}`$
$`(\stackrel{~}{L}_5^{k1,k,1}+2\stackrel{~}{L}_4^{k1,k,2}+2\stackrel{~}{L}_3^{k1,k,2}+\stackrel{~}{L}_2^{k1,k,1}),`$
$`V_1^{k1,k,4}(n;(1)+(1)+(1))`$ $`=`$ $`\stackrel{~}{L}_n^{k1,k,1}+3\stackrel{~}{L}_{n1}^{k1,k,1}+3\stackrel{~}{L}_{n2}^{k1,k,1}+\stackrel{~}{L}_{n3}^{k1,k,1}`$
$`(\stackrel{~}{L}_5^{k1,k,1}+3\stackrel{~}{L}_4^{k1,k,2}+3\stackrel{~}{L}_3^{k1,k,2}+\stackrel{~}{L}_2^{k1,k,1}),`$
$`V_2^{k1,k,4}(n;(1)+(1))`$ $`=`$ $`V_2^{k1,k,3}(n;(1))+V_2^{k1,k,3}(n1;(1))(\stackrel{~}{L}_5^{k1,k,2}\stackrel{~}{L}_3^{k1,k,2})`$ (4.53)
$`+{\displaystyle \frac{1}{2}}(V_1^{k1,k,1}(n;(1))(\stackrel{~}{L}_{n3}^{k1,k,1}\stackrel{~}{L}_2^{k1,k,1})`$
$`+(\stackrel{~}{L}_n^{k1,k,1}\stackrel{~}{L}_2^{k1,k,1})V_1^{k1,k,1}(n2;(1))`$
$`V_1^{k1,k,4}(n;(1)+(2))(\stackrel{~}{L}_3^{k1,k,1}\stackrel{~}{L}_2^{k1,k,1})`$
$`V_1^{k1,k,4}(n;(3))(\stackrel{~}{L}_4^{k1,k,1}\stackrel{~}{L}_2^{k1,k,1})).`$
###### Conjecture 3
In the $`d=4`$ case, $`V_{4m}^{k1,k,4}(n;\sigma _m)`$ equals $`G_{4m}^{k1,k,4}(n;\sigma _m)`$ except for $`V_2^{k1,k,4}(n,(1)+(1))`$. $`G_2^{k1,k,4}(n,(1)+(1))`$ is given by the formula:
$`G_2^{k1,k,4}(n;(1)+(1))`$ $`=`$ $`V_2^{k1,k,3}(n;(1))+V_2^{k1,k,3}(n1;(1))(\stackrel{~}{L}_5^{k1,k,2}\stackrel{~}{L}_3^{k1,k,2})`$ (4.54)
$`+{\displaystyle \frac{3}{4}}(V_1^{k1,k,1}(n;(1))(\stackrel{~}{L}_{n3}^{k1,k,1}\stackrel{~}{L}_2^{k1,k,1})`$
$`+(\stackrel{~}{L}_n^{k1,k,1}\stackrel{~}{L}_2^{k1,k,1})V_1^{k1,k,1}(n2;(1))`$
$`V_1^{k1,k,4}(n;(1)+(2))(\stackrel{~}{L}_3^{k1,k,1}\stackrel{~}{L}_2^{k1,k,1})`$
$`V_1^{k1,k,4}(n;(3))(\stackrel{~}{L}_4^{k1,k,1}\stackrel{~}{L}_2^{k1,k,1})).`$
The modification of the factor $`\frac{1}{2}`$ in (4.53) into the factor $`\frac{3}{4}`$ in (4.54) is determined by one numerical data:
$$L_7^{11,12,4}=1324882975682876246483412831870565329165165953902032,$$
(4.55)
and the generalized mirror transformation obtained by Conjecture 3 and (1.1) correctly predicts $`L_n^{k1,k,4}`$ up to $`k18`$.
In the $`d=5`$ case, we can state the following proposition under the assumption of Conjecture 2.
###### Proposition 7
$`G_5^{k1,k,5}(n;(0))=V_5^{k1,k,5}(n;(0)),`$
$`G_4^{k1,k,5}(n;(1))=V_4^{k1,k,5}(n;(1)),`$
$`G_3^{k1,k,5}(n;(2))=V_3^{k1,k,5}(n;(2)),`$
$`G_2^{k1,k,5}(n;(3))=V_2^{k1,k,5}(n;(3)),`$
$`G_1^{k1,k,5}(n;(4))=V_1^{k1,k,5}(n;(4)),`$
$`G_1^{k1,k,5}(n;(3)+(1))=V_1^{k1,k,5}(n;(3)+(1)),`$
$`G_1^{k1,k,5}(n;(2)+(2))=V_1^{k1,k,5}(n;(2)+(2)),`$
$`G_1^{k1,k,5}(n;(2)+(1)+(1))=V_1^{k1,k,5}(n;(2)+(1)+(1)),`$
$`G_1^{k1,k,5}(n;(1)+(1)+(1)+(1))=V_1^{k1,k,5}(n;(1)+(1)+(1)+(1)).`$ (4.56)
On the other hand, we find that there exist some non-trivial modifications to obtain $`G_3^{k1,k,5}(n;(1)+(1))`$, $`G_2^{k1,k,5}(n;(1)+(2))`$ and $`G_2^{k1,k,5}(n;(1)+(1)+(1))`$ from the corresponding virtual Gromov- Witten Invariants.
###### Conjecture 4
$`G_2^{k1,k,5}(n;(1)+(2))`$ $`=`$ $`V_2^{k1,k,4}(n;(2))+V_2^{k1,k,4}(n1;(2))\stackrel{~}{L}_6^{k1,k,2}+\stackrel{~}{L}_3^{k1,k,2}`$ (4.57)
$`+{\displaystyle \frac{8}{5}}hi_1(n)+hi_2(n)+{\displaystyle \frac{4}{5}}hi_3(n){\displaystyle \frac{3}{5}}hi_4(n),`$
$`G_2^{k1,k,5}(n;(1)+(1)+(1))`$ $`=`$ $`V_2^{k1,k,3}(n;(1))+V_2^{k1,k,3}(n1;(1))(\stackrel{~}{L}_5^{k1,k,2}\stackrel{~}{L}_3^{k1,k,2})`$ (4.58)
$`+{\displaystyle \frac{4}{5}}(V_1^{k1,k,1}(n;(1))(\stackrel{~}{L}_{n3}^{k1,k,1}\stackrel{~}{L}_2^{k1,k,1})`$
$`+(\stackrel{~}{L}_n^{k1,k,1}\stackrel{~}{L}_2^{k1,k,1})V_1^{k1,k,1}(n2;(1))`$
$`V_1^{k1,k,4}(n;(1)+(2))(\stackrel{~}{L}_3^{k1,k,1}\stackrel{~}{L}_2^{k1,k,1})`$
$`V_1^{k1,k,4}(n;(3))(\stackrel{~}{L}_4^{k1,k,1}\stackrel{~}{L}_2^{k1,k,1}))`$
$`+V_2^{k1,k,3}(n1;(1))+V_2^{k1,k,3}(n2;(1))(\stackrel{~}{L}_5^{k1,k,2}\stackrel{~}{L}_3^{k1,k,2})`$
$`+{\displaystyle \frac{4}{5}}(V_1^{k1,k,1}(n1;(1))(\stackrel{~}{L}_{n4}^{k1,k,1}\stackrel{~}{L}_2^{k1,k,1})`$
$`+(\stackrel{~}{L}_{n1}^{k1,k,1}\stackrel{~}{L}_2^{k1,k,1})V_1^{k1,k,1}(n3;(1))`$
$`V_1^{k1,k,4}(n1;(1)+(2))(\stackrel{~}{L}_3^{k1,k,1}\stackrel{~}{L}_2^{k1,k,1})`$
$`V_1^{k1,k,4}(n1;(3))(\stackrel{~}{L}_4^{k1,k,1}\stackrel{~}{L}_2^{k1,k,1}))V_2^{k1,k,3}(6;(1))`$
$`+{\displaystyle \frac{46}{25}}hi_1(n)+{\displaystyle \frac{46}{25}}hi_2(n)+{\displaystyle \frac{16}{25}}hi_3(n){\displaystyle \frac{2}{25}}hi_4(n),`$
$`G_3^{k1,k,5}(n;(1)+(1))`$ $`=`$ $`V_3^{k1,k,4}(n;(1))+V_3^{k1,k,4}(n1;(1))(\stackrel{~}{L}_6^{k1,k,3}\stackrel{~}{L}_4^{k1,k,3})`$
$`+{\displaystyle \frac{4}{5}}(V_2^{k1,k,3}(n;(1))(\stackrel{~}{L}_{n4}^{k1,k,1}\stackrel{~}{L}_2^{k1,k,1})`$
$`+(\stackrel{~}{L}_n^{k1,k,1}\stackrel{~}{L}_2^{k1,k,1})V_2^{k1,k,3}(n2;(1))`$
$`(V_2^{k1,k,4}(n;(2))+V_2^{k1,k,4}(n1;(2))`$
$`(\stackrel{~}{L}_6^{k1,k,2}\stackrel{~}{L}_3^{k1,k,2}))(\stackrel{~}{L}_3^{k1,k,1}\stackrel{~}{L}_2^{k1,k,1})`$
$`V_1^{k1,k,5}(n;(4))(\stackrel{~}{L}_5^{k1,k,2}\stackrel{~}{L}_3^{k1,k,2}))`$
$`+{\displaystyle \frac{3}{5}}(V_1^{k1,k,2}(n;(1))(\stackrel{~}{L}_{n3}^{k1,k,2}\stackrel{~}{L}_3^{k1,k,2})`$
$`+(\stackrel{~}{L}_n^{k1,k,2}\stackrel{~}{L}_3^{k1,k,2})V_1^{k1,k,2}(n3;(1))`$
$`V_1^{k1,k,5}(n;(1)+(3))(\stackrel{~}{L}_4^{k1,k,2}\stackrel{~}{L}_3^{k1,k,2})`$
$`V_2^{k1,k,5}(n;(3))(\stackrel{~}{L}_4^{k1,k,1}\stackrel{~}{L}_2^{k1,k,1}))`$
$`(\stackrel{~}{L}_3^{k1,k,1}\stackrel{~}{L}_2^{k1,k,1})({\displaystyle \frac{6}{5}}hi_1(n)+hi_2(n)+{\displaystyle \frac{3}{5}}hi_3(n){\displaystyle \frac{1}{5}}hi_4(n)),`$
where $`hi_j(n)`$ is a degree $`2`$ homogeneous polynomial of $`\stackrel{~}{L}_m^{k1,k,1}`$ satisfying $`hi_j(6)=hi_j(7)=0`$ and is given by,
$`hi_1(n)`$ $`=`$ $`\stackrel{~}{L}_n^{k1,k,1}\stackrel{~}{L}_{n4}^{k1,k,1}\stackrel{~}{L}_3^{k1,k,1}(\stackrel{~}{L}_n^{k1,k,1}+\stackrel{~}{L}_{n1}^{k1,k,1}+\stackrel{~}{L}_{n2}^{k1,k,1}+\stackrel{~}{L}_{n3}^{k1,k,1}+\stackrel{~}{L}_{n4}^{k1,k,1})`$
$`+\stackrel{~}{L}_2^{k1,k,1}(\stackrel{~}{L}_{n1}^{k1,k,1}+\stackrel{~}{L}_{n2}^{k1,k,1}+\stackrel{~}{L}_{n3}^{k1,k,1})`$
$`(\stackrel{~}{L}_6^{k1,k,1}\stackrel{~}{L}_2^{k1,k,1}\stackrel{~}{L}_3^{k1,k,1}(\stackrel{~}{L}_6^{k1,k,1}+\stackrel{~}{L}_5^{k1,k,1}+\stackrel{~}{L}_4^{k1,k,1}+\stackrel{~}{L}_3^{k1,k,1}+\stackrel{~}{L}_2^{k1,k,1})`$
$`+\stackrel{~}{L}_2^{k1,k,1}(\stackrel{~}{L}_5^{k1,k,1}+\stackrel{~}{L}_4^{k1,k,1}+\stackrel{~}{L}_3^{k1,k,1})),`$
$`hi_2(n)`$ $`=`$ $`\stackrel{~}{L}_n^{k1,k,1}\stackrel{~}{L}_{n3}^{k1,k,1}+\stackrel{~}{L}_{n1}^{k1,k,1}\stackrel{~}{L}_{n4}^{k1,k,1}`$
$`\stackrel{~}{L}_4^{k1,k,1}(\stackrel{~}{L}_n^{k1,k,1}+\stackrel{~}{L}_{n1}^{k1,k,1}+\stackrel{~}{L}_{n2}^{k1,k,1}+\stackrel{~}{L}_{n3}^{k1,k,1}+\stackrel{~}{L}_{n4}^{k1,k,1})+\stackrel{~}{L}_2^{k1,k,1}\stackrel{~}{L}_{n2}^{k1,k,1}`$
$`(\stackrel{~}{L}_6^{k1,k,1}\stackrel{~}{L}_3^{k1,k,1}+\stackrel{~}{L}_5^{k1,k,1}\stackrel{~}{L}_2^{k1,k,1}`$
$`\stackrel{~}{L}_4^{k1,k,1}(\stackrel{~}{L}_6^{k1,k,1}+\stackrel{~}{L}_5^{k1,k,1}+\stackrel{~}{L}_4^{k1,k,1}+\stackrel{~}{L}_3^{k1,k,1}+\stackrel{~}{L}_2^{k1,k,1})+\stackrel{~}{L}_2^{k1,k,1}\stackrel{~}{L}_4^{k1,k,1}),`$
$`hi_3(n)`$ $`=`$ $`\stackrel{~}{L}_n^{k1,k,1}\stackrel{~}{L}_{n2}^{k1,k,1}+\stackrel{~}{L}_{n1}^{k1,k,1}\stackrel{~}{L}_{n3}^{k1,k,1}+\stackrel{~}{L}_{n2}^{k1,k,1}\stackrel{~}{L}_{n4}^{k1,k,1}`$
$`\stackrel{~}{L}_5^{k1,k,1}(\stackrel{~}{L}_n^{k1,k,1}+\stackrel{~}{L}_{n1}^{k1,k,1}+\stackrel{~}{L}_{n2}^{k1,k,1}+\stackrel{~}{L}_{n3}^{k1,k,1}+\stackrel{~}{L}_{n4}^{k1,k,1})\stackrel{~}{L}_2^{k1,k,1}\stackrel{~}{L}_{n2}^{k1,k,1}`$
$`(\stackrel{~}{L}_6^{k1,k,1}\stackrel{~}{L}_4^{k1,k,1}+\stackrel{~}{L}_5^{k1,k,1}\stackrel{~}{L}_3^{k1,k,1}+\stackrel{~}{L}_4^{k1,k,1}\stackrel{~}{L}_2^{k1,k,1}`$
$`\stackrel{~}{L}_5^{k1,k,1}(\stackrel{~}{L}_6^{k1,k,1}+\stackrel{~}{L}_5^{k1,k,1}+\stackrel{~}{L}_4^{k1,k,1}+\stackrel{~}{L}_3^{k1,k,1}+\stackrel{~}{L}_2^{k1,k,1})\stackrel{~}{L}_2^{k1,k,1}\stackrel{~}{L}_4^{k1,k,1}),`$
$`hi_4(n)`$ $`=`$ $`\stackrel{~}{L}_{n1}^{k1,k,1}\stackrel{~}{L}_{n3}^{k1,k,1}\stackrel{~}{L}_4^{k1,k,1}(\stackrel{~}{L}_{n1}^{k1,k,1}+\stackrel{~}{L}_{n2}^{k1,k,1}+\stackrel{~}{L}_{n3}^{k1,k,1})+\stackrel{~}{L}_3^{k1,k,1}\stackrel{~}{L}_{n2}^{k1,k,1}`$
$`(\stackrel{~}{L}_5^{k1,k,1}\stackrel{~}{L}_3^{k1,k,1}\stackrel{~}{L}_4^{k1,k,1}(\stackrel{~}{L}_5^{k1,k,1}+\stackrel{~}{L}_4^{k1,k,1}+\stackrel{~}{L}_3^{k1,k,1})+\stackrel{~}{L}_3^{k1,k,1}\stackrel{~}{L}_4^{k1,k,1}).`$
As in the $`d=4`$ case, we have fixed the modification of the rational factors by one numerical data:
$`L_8^{12,13,5}=`$
$`100355724573836807695163109854598526931747042477505803923089934593470758513921/180000.`$
Check of the prediction formula obtained by (1.1), Proposition 7 and Conjecture 4 takes a lot of time due to numerical computation of fixed point formulas. We checked that the formula correctly predicts $`L_8^{13,14,5}`$. Now, we discuss the rules of modification of $`V_{dm}^{k1,k,d}(n;\sigma _m)`$ into $`G_{dm}^{k1,k,d}(n;\sigma _m)`$. First, the corresponding $`V_{5m}^{k1,k,5}(n;\sigma _m)`$ is given as follows.
$`V_2^{k1,k,5}(n;(1)+(2))`$ $`=`$ $`V_2^{k1,k,4}(n;(2))+V_2^{k1,k,4}(n1;(2))\stackrel{~}{L}_6^{k1,k,2}+\stackrel{~}{L}_3^{k1,k,2}`$ (4.62)
$`+{\displaystyle \frac{1}{2}}(2hi_1(n)+hi_2(n)+hi_3(n)hi_4(n)),`$
$`V_2^{k1,k,5}(n;(1)+(1)+(1))`$ $`=`$ $`V_2^{k1,k,3}(n;(1))+V_2^{k1,k,3}(n1;(1))(\stackrel{~}{L}_5^{k1,k,2}\stackrel{~}{L}_3^{k1,k,2})`$ (4.63)
$`+{\displaystyle \frac{1}{2}}(V_1^{k1,k,1}(n;(1))(\stackrel{~}{L}_{n3}^{k1,k,1}\stackrel{~}{L}_2^{k1,k,1})`$
$`+(\stackrel{~}{L}_n^{k1,k,1}\stackrel{~}{L}_2^{k1,k,1})V_1^{k1,k,1}(n2;(1))`$
$`V_1^{k1,k,4}(n;(1)+(2))(\stackrel{~}{L}_3^{k1,k,1}\stackrel{~}{L}_2^{k1,k,1})`$
$`V_1^{k1,k,4}(n;(3))(\stackrel{~}{L}_4^{k1,k,1}\stackrel{~}{L}_2^{k1,k,1}))`$
$`+V_2^{k1,k,3}(n1;(1))+V_2^{k1,k,3}(n2;(1))(\stackrel{~}{L}_5^{k1,k,2}\stackrel{~}{L}_3^{k1,k,2})`$
$`+{\displaystyle \frac{1}{2}}(V_1^{k1,k,1}(n1;(1))(\stackrel{~}{L}_{n4}^{k1,k,1}\stackrel{~}{L}_2^{k1,k,1})`$
$`+(\stackrel{~}{L}_{n1}^{k1,k,1}\stackrel{~}{L}_2^{k1,k,1})V_1^{k1,k,1}(n3;(1))`$
$`V_1^{k1,k,4}(n1;(1)+(2))(\stackrel{~}{L}_3^{k1,k,1}\stackrel{~}{L}_2^{k1,k,1})`$
$`V_1^{k1,k,4}(n1;(3))(\stackrel{~}{L}_4^{k1,k,1}\stackrel{~}{L}_2^{k1,k,1}))V_2^{k1,k,3}(6;(1))`$
$`+{\displaystyle \frac{1}{4}}(3hi_1(n)+3hi_2(n)+hi_3(n)),`$
$`V_3^{k1,k,5}(n;(1)+(1))`$ $`=`$ $`V_3^{k1,k,4}(n;(1))+V_3^{k1,k,4}(n1;(1))(\stackrel{~}{L}_6^{k1,k,3}\stackrel{~}{L}_4^{k1,k,3})`$
$`+{\displaystyle \frac{2}{3}}(V_2^{k1,k,3}(n;(1))(\stackrel{~}{L}_{n4}^{k1,k,1}\stackrel{~}{L}_2^{k1,k,1})`$
$`+(\stackrel{~}{L}_n^{k1,k,1}\stackrel{~}{L}_2^{k1,k,1})V_2^{k1,k,3}(n2;(1))`$
$`(V_2^{k1,k,4}(n;(2))+V_2^{k1,k,4}(n1;(2))`$
$`(\stackrel{~}{L}_6^{k1,k,2}\stackrel{~}{L}_3^{k1,k,2}))(\stackrel{~}{L}_3^{k1,k,1}\stackrel{~}{L}_2^{k1,k,1})`$
$`V_1^{k1,k,5}(n;(4))(\stackrel{~}{L}_5^{k1,k,2}\stackrel{~}{L}_3^{k1,k,2}))`$
$`+{\displaystyle \frac{1}{3}}(V_1^{k1,k,2}(n;(1))(\stackrel{~}{L}_{n3}^{k1,k,2}\stackrel{~}{L}_3^{k1,k,2})`$
$`+(\stackrel{~}{L}_n^{k1,k,2}\stackrel{~}{L}_3^{k1,k,2})V_1^{k1,k,2}(n3;(1))`$
$`V_1^{k1,k,5}(n;(1)+(3))(\stackrel{~}{L}_4^{k1,k,2}\stackrel{~}{L}_3^{k1,k,2})`$
$`V_2^{k1,k,5}(n;(3))(\stackrel{~}{L}_4^{k1,k,1}\stackrel{~}{L}_2^{k1,k,1}))`$
$`{\displaystyle \frac{1}{3}}(\stackrel{~}{L}_3^{k1,k,1}\stackrel{~}{L}_2^{k1,k,1})(2hi_1(n)+hi_2(n)+hi_3(n)hi_4(n))`$
The non-trivial rational coefficients of $`V_{dm}^{k1,k,d}(n;\sigma _m)`$ come from the rational factor,
$$\underset{j=1}{\overset{l(\sigma _m)1}{}}(1\frac{n_j}{dm}),(0n_j<dm),$$
(4.65)
whose origin is the factor $`\frac{1}{(dm)^{l(\sigma _m)1}}`$ in the definition of $`V_{dm}^{k1,k,d}(n;\sigma _m)`$. Then looking at the above formulas of $`V_{5m}^{k1,k,5}(n;\sigma _m)`$ and $`G_{5m}^{k1,k,5}(n;\sigma _m)`$, we can speculate that the rational factor in (4.65) is modified into,
$$\underset{j=1}{\overset{l(\sigma _m)1}{}}(1\frac{n_j}{d}).$$
(4.66)
Let us take $`V_2^{k1,k,5}(n;(1)+(1)+(1))`$ as an example. We decompose $`V_2^{k1,k,5}(n;(1)+(1)+(1))`$ according to the decomposition in (3.44):
$`V_2^{k1,k,5}(n;(1)+(1)+(1))`$ $`=`$ $`\pi _1\pi _1\pi _1(\stackrel{~}{L}_n^{k1,k,2}\stackrel{~}{L}_3^{k1,k,2})+\pi _1\pi _1(hi_2^{k1,k,3}(n;(1)))`$
$`+\pi _1(hi_2^{k1,k,4}(n;(1)+(1)))+hi_2^{k1,k,5}(n;(1)+(1)+(1)).`$
More explicitly, each decomposed part corresponds to,
$`\pi _1\pi _1\pi _1(\stackrel{~}{L}_n^{k1,k,2}\stackrel{~}{L}_3^{k1,k,2})+\pi _1\pi _1(hi_2^{k1,k,3}(n;(1)))=`$
$`V_2^{k1,k,3}(n;(1))+V_2^{k1,k,3}(n1;(1))(\stackrel{~}{L}_5^{k1,k,2}\stackrel{~}{L}_3^{k1,k,2})`$
$`+V_2^{k1,k,3}(n1;(1))+V_2^{k1,k,3}(n2;(1))(\stackrel{~}{L}_5^{k1,k,2}\stackrel{~}{L}_3^{k1,k,2})V_2^{k1,k,3}(6;(1)),`$
$`\pi _1(hi_2^{k1,k,4}(n;(1)+(1)))=`$
$`{\displaystyle \frac{1}{2}}(V_1^{k1,k,1}(n;(1))(\stackrel{~}{L}_{n3}^{k1,k,1}\stackrel{~}{L}_2^{k1,k,1})+(\stackrel{~}{L}_n^{k1,k,1}\stackrel{~}{L}_2^{k1,k,1})V_1^{k1,k,1}(n2;(1))`$
$`V_1^{k1,k,4}(n;(1)+(2))(\stackrel{~}{L}_3^{k1,k,1}\stackrel{~}{L}_2^{k1,k,1})V_1^{k1,k,4}(n;(3))(\stackrel{~}{L}_4^{k1,k,1}\stackrel{~}{L}_2^{k1,k,1})`$
$`+V_1^{k1,k,1}(n1;(1))(\stackrel{~}{L}_{n4}^{k1,k,1}\stackrel{~}{L}_2^{k1,k,1})+(\stackrel{~}{L}_{n1}^{k1,k,1}\stackrel{~}{L}_2^{k1,k,1})V_1^{k1,k,1}(n3;(1))`$
$`V_1^{k1,k,4}(n1;(1)+(2))(\stackrel{~}{L}_3^{k1,k,1}\stackrel{~}{L}_2^{k1,k,1})V_1^{k1,k,4}(n1;(3))(\stackrel{~}{L}_4^{k1,k,1}\stackrel{~}{L}_2^{k1,k,1})),`$
$`hi_2^{k1,k,5}(n;(1)+(1)+(1))={\displaystyle \frac{1}{4}}(3hi_1(n)+3hi_2(n)+hi_3(n)).`$ (4.70)
According to (4.66), the modification is given by,
$$\frac{1}{2}\frac{4}{5},\frac{1}{4}=(\frac{1}{2})^2(\frac{4}{5})^2=\frac{16}{25}.$$
(4.71)
This modification is almost correct, but there exist some errors in the modification of (4.70), which are given by,
$$\frac{2}{25}(hi_1(n)+hi_2(n)+hi_4(n)).$$
(4.72)
Similar errors also occur in the cases of $`V_2^{k1,k,5}(n;(1)+(2))`$ and $`V_3^{k1,k,5}(n;(1)+(1))`$. One of the reasons of such errors comes from the fact that the decomposition of (3.44) is not unique, as was suggested in the remark of (3.44), but there must be other reasons because in the case of $`V_2^{k1,k,5}(n;(1)+(1)+(1))`$, the decomposition is unique. Therefore, further consideration is needed.
###### Question 1
Fix the general rule of the modification of $`V_{dm}^{N,k,d}(n;\sigma _m)`$ into $`G_{dm}^{N,k,d}(n;\sigma _m)`$.
###### Remark 4
$`\pi _1(hi_2^{k1,k,4}(n;(1)+(1)))`$ does not vanish when $`n=7`$, but Conjecture 4 tells us that it receives modification. Thus, the condition (iv) of Proposition 1 does not hold true for $`G_{dm}^{N,k,d}(n;\sigma _m)`$ in the $`d5`$ cases.
Acknowledgement
We would like to thank T.Eguchi, S.Hosono and A.Collino for discussions. We also thank M.Naka for kind encouragement. Research of the author is supported by JSPS postodoctoral fellowship. |
warning/0002/math0002036.html | ar5iv | text | # Wave invariants at elliptic closed geodesics
## 0. Introduction
The purpose of this article is to provide an effective method for calculating the wave invariants associated to a non-degenerate elliptic closed geodesic $`\gamma `$ of a compact Riemannian manifold $`(M,g)`$: that is, the coefficients in the singularity expansion
$$TrU(t)=e_o(t)+\underset{\gamma }{}e_\gamma (t)$$
$$e_\gamma (t)a_{\gamma 1}(tL_\gamma +i0)^1+\underset{k=0}{\overset{\mathrm{}}{}}a_{\gamma k}(tL_\gamma +i0)^klog(tL_\gamma +i0)$$
of the trace of the wave group $`U(t):=expit\sqrt{\mathrm{\Delta }}`$ at $`t=L_\gamma `$ (the length of $`\gamma .`$) We will show that
$$a_{\gamma k}=_\gamma I_{\gamma ;k}(s)𝑑s$$
for certain homogeneous invariant densities $`I_{\gamma k}(s)ds`$ on $`\gamma `$, given by at most 2k+1 integrals over $`\gamma `$ of polynomials in the curvature, Jacobi fields, length, inverse length and Floquet invariants $`\beta _j:=(1e^{i\alpha _j})^1`$ along $`\gamma .`$ These expressions characterize the wave invariants in much the same way that the heat invariants (or wave invariants at t=0) are characterized as integrals $`_MP_j(R,R,\mathrm{})𝑑vol`$ of homogeneous curvature polynomials over M \[ABP\] \[Gi\]. Moreover, in combination with the recent inverse results of Guillemin \[G.1,2\], the method produces a list of new spectral invariants of this kind, simpler than the wave invariants themselves (the so-called quantum Birkhoff normal form coefficients $`B_{\gamma k;j}`$.)
To state the results, we will need some notation. We let $`𝒥_\gamma ^{}`$ denote the space of complex normal Jacobi fields along $`\gamma `$, a symplectic vector space of (complex) dimension 2n (n=dim M-1) with respect to the Wronskian
$$\omega (X,Y)=g(X,\frac{D}{ds}Y)g(\frac{D}{ds}X,Y).$$
The linear Poincare map $`P_\gamma `$ is then the linear symplectic map on $`𝒥_\gamma ^{}`$ defined by $`P_\gamma Y(t)=Y(t+L_\gamma ).`$ We will assume $`\gamma `$ to be non-degenerate elliptic, i.e. that the eigenvalues of $`P_\gamma `$ are of the form $`\{e^{\pm i\alpha _j},j=1,\mathrm{},n\}`$ with (Floquet) exponents $`\{\alpha _1,\mathrm{},\alpha _n\}`$, together with $`\pi `$, independent over $`𝐐`$. The associated normalized eigenvectors will be denoted $`\{Y_j,\overline{Y_j},j=1,\mathrm{},n\}`$,
$$P\gamma Y_j=e^{i\alpha _j}Y_jP_\gamma \overline{Y}_j=e^{i\alpha _j}\overline{Y}_j\omega (Y_j,\overline{Y}_k)=\delta _{jk}$$
and relative to a fixed parallel normal frame $`e(s):=(e_1(s),\mathrm{},e_n(s))`$ along $`\gamma `$ they will be written in the form $`Y_j(s)=_{k=1}^ny_{jk}(s)e_k(s).`$ The metric coefficients $`g_{ij}`$ will always be taken relative to Fermi normal coordinates $`(s,y)`$ along $`\gamma `$. The mth jet of $`g`$ along $`\gamma `$ will be denoted by $`j_\gamma ^mg`$, the curvature tensor by $`R`$ and its covariant derivatives by $`^mR`$. The vector fields $`\frac{}{s},\frac{}{y_j}`$ and their real linear combinations will be referred to as Fermi normal vector fields along $`\gamma `$ and contractions of tensor products of the $`^mR`$’s with these vector fields will be referred to as Fermi curvature polynomials. Such polynomials will be called invariant if they are invariant under the action of $`O(n)`$ in the normal spaces. Invariant contractions against $`\frac{}{s}`$ and against the Jacobi eigenfields $`Y_j,\overline{Y}_j`$, with coefficients given by invariant polynomials in the components $`y_{jk}`$, will be called Fermi-Jacobi polynomials. We will alse use this term for functions on $`\gamma `$ given by repeated indefinite integrals over $`\gamma `$ of such FJ polynomials. Finally, FJ polynomials whose coefficients are given by polynomials in the Floquet invariants $`\beta _j=(1e^{i\alpha _j})^1`$ will be called Fermi-Jacobi-Floquet polynomials. We give weights to the variables $`g_{ij},D_{s,y}^\beta g_{ij},L:=L_\gamma ,\alpha _j,y_{ij},\dot{y}_{ij}`$ as follows: $`wgt(D_{s,y}^\beta g_{ij})=|\beta |,wgt(L)=1,wgt(\alpha _j)=0,wgt(y_{ij})=\frac{1}{2},wgt(\dot{y}_{ij})=\frac{1}{2}.`$ As will be seen, these weights reflect the scaling of these objects under $`gϵ^2g`$. A polynomial in this data is homogeneous of weight s if all its monomials have weight s under this scaling.
Theorem A Let $`\gamma `$ be an elliptic closed geodesic with $`\{\alpha _1,\mathrm{},\alpha _n,\pi \}`$ independent over $`𝐐`$. Then $`a_{\gamma k}=_\gamma I_{\gamma k}(s;g)𝑑s`$ where :
(i) $`I_{\gamma k}(s,g)`$ is a homogeneous Fermi-Jacobi-Floquet polynomial of weight -k - 1 in the data $`\{L,y_{ij},\dot{y}_{ij},D_{s,y}^\beta g\}`$ with $`|\beta |2k+4`$ ;
(ii) The degree of $`I_{\gamma k}`$ in the Jacobi field components is at most 6k+6;
(iii) At most 2k+1 indefinite integrations over $`\gamma `$ occur in $`I_{\gamma k}`$;
(iv) The degree of $`I_{\gamma k}`$ in the Floquet invariants $`\beta _j`$ is at most k+2.
For instance, in dimension 2 the residual wave invariant $`a_{\gamma o}`$ is given by:
$$a_{\gamma o}=\frac{c_\gamma }{L^\mathrm{\#}}[B_{\gamma o;4}(2\beta ^2\beta \frac{3}{4})+B_{\gamma o;0}]$$
where:
(a) $`c_\gamma `$ is the principal wave invariant $`i^\sigma L^\mathrm{\#}|IP_\gamma |^{\frac{1}{2}}`$;
(b) $`L^\mathrm{\#}`$ the primitive length of $`\gamma `$; $`\sigma `$ is its Morse index; $`P_\gamma `$ is its Poincare map;
(c) $`B_{\gamma o;j}`$ have the form:
$$B_{\gamma o;j}=\frac{1}{L^\mathrm{\#}}_o^{L^\mathrm{\#}}[a|\dot{Y}|^4+b_1\tau |\dot{Y}Y|^2+b_2\tau Re(\overline{Y}\dot{Y})^2+c\tau ^2|Y|^4+d\tau _{\nu \nu }|Y|^4+e\delta _{jo}\tau ]𝑑s$$
$$+\frac{1}{L^\mathrm{\#}}\underset{0m,n3;m+n=3}{}C_{1;mn}\frac{sin((nm)\alpha )}{|(1e^{i(mn)\alpha })|^2}|_o^{L^\mathrm{\#}}\tau _\nu (s)\overline{Y}^mY^n(s)𝑑s|^2$$
$$+\frac{1}{L^\mathrm{\#}}\underset{0m,n3;m+n=3}{}C_{2;mn}Im\{_o^{L^\mathrm{\#}}\tau _\nu (s)\overline{Y}^mY^n(s)[_o^s\tau _\nu (t)\overline{Y}^nY^m(t)𝑑t]𝑑s\}$$
for various universal (computable) coefficients. Here,
(d) $`\tau `$ denotes the scalar curvature, $`\tau _\nu `$ its unit normal derivative, $`\tau _{\nu \nu }`$ the Hessian $`Hess(\tau )(\nu ,\nu )`$; $`Y`$ denotes the unique normalized Jacobi eigenfield, $`\dot{Y}`$ its time-derivative and $`\delta _{jo}`$ the Kronecker symbol (1 if j=0 and otherwise 0.)
We note that the residual wave invariant already saturates the description in Theorem A.
This characterization of the wave invariants makes more concrete (for Laplacians) the recent results of Guillemin \[G1,2\] which show that the wave invariants may be expressed in terms of the quantum Birkhoff normal form coefficients of the wave operator around $`\gamma .`$ For instance, Guillemin shows \[G.2, (8.24)\]:
$$a_{\gamma o}=ic_\gamma \sigma [\underset{ij}{}\frac{^2}{I_iI_j}H_1(0,\sigma )(\beta _i+\frac{1}{2})(\beta _j+\frac{1}{2})$$
$$+\frac{1}{2}\underset{i}{}(2\beta _i^2+\beta _i\frac{1}{4})\frac{^2}{I_i^2}H_1(0,\sigma )\underset{i}{}(\beta _i+\frac{1}{2})\frac{}{I_i}H_o(0,\sigma )+H_1(0,\sigma )]$$
where $`H_{1r}(\sigma ,I_1,\mathrm{},I_n)`$ is the term of order 1-r in the complete symbol of the quantum normal form. The coefficients $`B_{\gamma o;j}`$ above are essentially these QBNF (quantum Birkhoff normal form) coefficients.
Indeed, the first step in the proof of Theorem A (which we call Theorem B) consists in explicitly constructing the QBNF for $`\sqrt{\mathrm{\Delta }}`$ near $`\gamma `$. The method is different from that in \[G.1,2\] and effectively calculates the QBNF coefficients as integrals of Fermi-Jacobi polynomials. It is based to some extent on the construction of a complete set of quasi-modes associated to $`\gamma `$ as presented in Babich-Buldyrev \[B.B\] although the emphasis is on intertwining operators rather than on quasi-modes per se. In the construction of the intertwining operators it also employs several ideas in Sjostrand \[Sj\], although it does not begin by putting the principal symbol in Birkhoff normal form as is done in \[Sj\] and \[G.2\].
The second step (which we call Theorem C) consists in calculating the wave invariants from the normal form. This is possible because the wave invariants are non-commutative residues of the wave operator and its $`t`$ derivatives and hence are invariant under conjugation by unitary Fourier Integral operators (cf. \[G.1,2\], \[Z.1\]). The residues of the normal form wave group will be easily seen to be polynomials in the QBNF coefficients and in the $`\beta _j`$’s. Some dimensional analysis of the wave coefficients then leads to the description in the statement of Theorem A.
Guillemin \[loc.cit\] has also proved the remarkable inverse result that, conversely, the QBNF coefficients can be determined from the wave invariants associated to $`\gamma `$ and its iterates, and therefore are themselves spectral invariants. In view of Theorems A-B this gives a list of new spectral invariants, for instance the QBNF coefficients $`B_{\gamma o;j}`$ above, in the form of geodesic integrals of FJ polynomials,
Let us now describe the ingredients of the proofs more precisely. Henceforth we will reserve the notation $`\gamma `$ for a primitive closed geodesic and will denote its iterates by $`\gamma ^m`$. To introduce the quantum Birkhoff normal form at $`\gamma `$ and its role in the calculations, we note first that the wave invariants associated to $`\gamma `$ are determined by the microlocalization of $`\mathrm{\Delta }`$ to a conic neighborhood
$`(0.1)`$
$$|y|<ϵ\frac{|\eta |}{\sigma }<ϵ$$
of the cone $`^+\gamma `$ thru $`\gamma `$ in $`T^{}M0`$. Here, $`(\sigma ,\eta )`$ denote the symplectically dual coordinates to the Fermi normal coordinates $`(s,y)`$ above. The microlocalization of $`\mathrm{\Delta }`$ to (0.1) is then given by
$$\mathrm{\Delta }_\psi :=\psi (s,D_s,y,D_y)^{}\mathrm{\Delta }\psi (s,D_s,y,D_y)$$
where $`D_{x_j}:=\frac{}{ix_j}`$ and where $`\psi (s,\sigma ,y,\eta )`$ is supported in (0.1) and identically one in some smaller conic neighborhoood . Often we omit explicit mention of the microlocal cut-off $`\psi `$ in calculations which are valid on its microsupport. Under the exponential map
$$exp:N_\gamma M$$
along the normal bundle to $`\gamma `$, the localization of $`g`$, resp. $`\mathrm{\Delta }`$, to the tubular neighborhood $`|y|<ϵ`$ pulls back to a locally well-defined metric, respectively Laplacian, on a similar neighborhood in $`N_\gamma `$. Hence $`exp`$ conjugates $`\mathrm{\Delta }_\psi `$ to an isometric microlocalized Laplacian on $`N_\gamma `$, which we continue to denote by $`\mathrm{\Delta }_\psi `$. We are thus reduced to calculating the wave invariants of a Laplacian on the model space $`S_L^1\times ^n`$ at the closed geodesic $`\gamma =S_L^1\times \{0\}`$, where $`S_L^1:=/L.`$ For the sake of simplicity we assume the normal bundle is orientable but note that the reduction is valid for immersed as well as embedded closed geodesics.
We now wish to put $`\mathrm{\Delta }_\psi `$ into normal form, which is first of all to conjugate it (modulo a small error) into a distinguished maximal abelian algebra $`𝒜`$ of pseudodifferential operators on the model space $`S_L^1\times ^n`$. Roughly speaking, $`𝒜`$ is generated by the tangential operator $`D_s:=\frac{}{is}`$ on $`S_L^1`$ together with the transverse harmonic oscillators
$`(0.2)`$
$$I_j=I_j(y,D_y):=\frac{1}{2}(D_{y_j}^2+y_j^2).$$
In the construction of the normal form, a special role will be played by the distinguished element
$`(0.3)`$
$$:=\frac{1}{L}(LD_s+H_\alpha )$$
where
$`(0.4)`$
$$H_\alpha :=\frac{1}{2}\underset{k=1}{\overset{n}{}}\alpha _kI_k$$
and where the choice of sign in $`\pm \alpha _k`$ will be specified below. This element comes up naturally as the semi-classical parameter in the construction of quasi-modes, although $`D_s`$ is more suitable for analysing the wave invariants. Note that both are elliptic elements in the conic neighborhood
$`(0.5)`$
$$|I_j|<ϵ\sigma I_j(y,\eta ):=\frac{1}{2}(y_j^2+\eta _j^2),$$
which will be the image of (0.1) under the conjugation to normal form.
The classical Birkhoff normal form theorem states roughly the following: near a non- degenerate elliptic closed geodesic $`\gamma `$ the Hamiltonian
$$H(x,\xi )=|\xi |:=\sqrt{\underset{ij=1}{\overset{n+1}{}}g^{ij}\xi _i\xi _j}$$
can be conjugated by a homogeneous local canonical transformation $`\chi `$ to the normal form
$`(0.6)`$
$$\chi ^{}H\sigma +\frac{1}{L}\underset{i,j=1}{\overset{n}{}}\alpha _jI_j+\frac{p_1(I_1,\mathrm{},I_n)}{\sigma }+\mathrm{}modO_{\mathrm{}}^1$$
where $`p_k`$ is homogeneous of order k+1 in $`I_1,\mathrm{},I_n`$, and where $`O_{\mathrm{}}^1`$ is the space of germs of functions homogeneous of degree 1 which vanish to infinite order along $`\gamma .`$ Note that all the terms in (0.6) are homogenous of degree 1 in $`(\sigma ,I_1,\mathrm{},I_n)`$, and that the order of vanishing at $`|I|=0`$ equals one plus the order of decay in $`\sigma `$. The coefficients of the monomials in the $`p_j(I_1,\mathrm{},I_n)`$ are known as the classical Birkhoff normal form invariants. (See the Appendix for some further details).
The quantum Birkhoff normal form is the more or less analogous statement on the operator level. In the following the symbol $``$ means that the two sides agree modulo operators whose complete symbols are of order 1 and vanish to infinite order on $`\gamma .`$ Also, $`O_m\mathrm{\Psi }^r`$ denotes the space of pseudodifferential operators of order $`r`$ whose complete symbols vanish to order $`m`$ at $`(y,\eta )=(0,0).`$
Theorem B There exists a microlocally elliptic Fourier Integral operator $`W`$ from the conic neighborhood (0.1) of $`^+\gamma `$ in $`T^{}N_\gamma 0`$ to the conic neighborhood (0.5) of $`T_+^{}S_L^1`$ in $`T^{}(S_L^1\times R^n)`$ such that:
$$𝒟:=W\sqrt{\mathrm{\Delta }_\psi }W^1\overline{\psi }(,I_1,\mathrm{},I_n)[+\frac{p_1(I_1,\mathrm{}.,I_n)}{L}+\frac{p_2(I_1,\mathrm{},I_n)}{(L)^2}+\mathrm{}+\frac{p_{k+1}(I_1,\mathrm{},I_n)}{(L)^{k+1}}+\mathrm{}]$$
$$D_s+\frac{1}{L}H_\alpha +\frac{\stackrel{~}{p}_1(I_1,\mathrm{},I_n)}{LD_s}+\frac{\stackrel{~}{p}_2(I_1,\mathrm{},I_n)}{(LD_s)^2}+\mathrm{}+\frac{\stackrel{~}{p}_{k+1}(I_1,\mathrm{},I_n)}{(LD_s)^{k+1}}+$$
where the numerators $`p_j(I_1,\mathrm{},I_n),\stackrel{~}{p}_j(I_1,\mathrm{},I_n)`$ are polynomials of degree j+1 in the variables $`I_1,\mathrm{},I_n`$, where $`\overline{\psi }`$ is microlocally supported in (0.5), and where $`W^1`$ denotes a microlocal inverse to $`W`$ in (0.5). The kth remainder term lies in the space $`_{j=o}^{k+2}O_{2(k+2j)}\mathrm{\Psi }^{1j}`$.
The QBNF coefficients will by definition be the coefficients of the monomials in the classical action ($`I_j`$-) variables in the complete Weyl symbols of the operators $`\stackrel{~}{p}_j(I_1,\mathrm{},I_n)`$. As mentioned above, the proof of Theorem B gives an effective method for calculating them as integrals over $`\gamma `$ of FJ polynomials.
The asymptotic relation in the above expansion may be viewed in either of two ways: First, as mentioned in the statement of the Theorem, the kth remainder is a sum of terms in $`𝒜`$ of orders $`1,0,\mathrm{},(k+1)`$ where the complete symbol of the term of order $`1j`$ must vanish to order $`2(k+2j)`$. This characterization of the remainder will play the key role in the calculation of the wave invariants, since terms in the normal form with low pseudodifferential order or with high vanishing order make no contribution to a given wave invariant. On the other hand, it may be viewed as a semi-classical asymptotic relation with $``$ playing the role of semi-classical parameter; thus the theorem gives a semi-classical expansion for $`𝒟`$ in terms of $``$. This point of view comes up naturally in the theory of quasi-modes associated to $`\gamma :`$ Indeed, consider the joint $`𝒜`$\- eigenfunctions
$`(0.7)`$
$$\varphi _{kq}^o(s,y):=e_k(s)\gamma _q(y)$$
with $`e_k(s):=e^{\frac{2\pi }{L}iks}`$ and with $`\gamma _q`$ $`qth`$ normalized Hermite function ($`q^n`$). The corresponding eigenvalues of $`𝒟`$ then have the semi-classical expansions
$`(0.8)`$
$$\lambda _{kq}r_{kq}+\frac{p_1(q)}{r_{kq}}+\frac{p_2(q)}{r_{kq}^2}+\mathrm{}$$
where
$`(0.9)`$
$$r_{kq}=\frac{1}{L}(2\pi k+\underset{j=1}{\overset{n}{}}(q_j+\frac{1}{2})\alpha _j)$$
are the eigenvalues of $``$. Here the index $`q`$ is held fixed as $`k\mathrm{}`$. We recognize in (0.8) the familiar form of the quasi-eigenvalues associated to $`\gamma `$ (cf. \[B.B., ch.9\]); hence the intertwining operator $`W`$ is the operator taking the eigenfunctions (0.7) to quasi-modes of infinite order for $`\mathrm{\Delta }`$ at $`\gamma .`$
As will be seen in (§4), Theorem B implies that the wave invariants of $`\sqrt{\mathrm{\Delta }}_\psi `$ are the same as the wave invariants of $`𝒟`$. The second main step in the calculation of the wave invariants is then the use of the non-commutative residue to connect the terms in the normal form expansion with the terms in the singularity expansion for $`Tr`$ exp $`it𝒟.`$ The main point here is that
$`(0.10)`$
$$a_{\gamma k}=resD_t^ke^{it\sqrt{\mathrm{\Delta }_\psi }}:=Res_{s=0}TrD_t^ke^{it\sqrt{\mathrm{\Delta }_\psi }}\sqrt{\mathrm{\Delta }_\psi }^s,$$
with $`res`$ invariant under conjugation by (microlocal) unitary operators, and depending on only a finite jet of the Laplacian near $`\gamma .`$ Hence it may be calculated by conjugating to the normal form, and indeed will only depend on a finite part of the normal form. Applying $`D_t^k`$ and formally exponentiating the terms of order $`1`$ in $`D_s`$ we get
$`(0.11)`$
$$res\overline{\psi }(D_s,I_1,\mathrm{},I_n)D_t^ke^{it𝒟}|_{t=L}=res\overline{\psi }(D_s,I_1,\mathrm{},I_n)e^{iLD_s}e^{iH_\alpha }𝒟^k(I+iL\frac{\stackrel{~}{p}_1(I_1,\mathrm{},I_n)}{LD_s}+\mathrm{})$$
which suggests that the wave coefficient $`a_{k\gamma }`$ is the regularized trace of the coefficient of $`D_{s}^{}{}_{}{}^{1}`$ in (0.11). This is not clear, even formally, since many of the terms of negative order in $`D_s`$ in the exponent have overall order 1 as pseudodifferential operators; but it will prove to be the case. Since $`e^{iLD_s}=I`$ on $`/L𝐙`$ the Fourier Integral factor in (0.11) is just $`e^{iH_\alpha }`$. Regarding the regularized traces of the coefficients, we note that
$`(0.12)`$
$$Tre^{iH_\alpha }=\underset{q𝐍^n}{}e^{i_{k=1}^n(q_k+\frac{1}{2})\alpha _k}$$
is well-defined as the tempered distribution
$`(0.13)`$
$$T(\alpha )=\mathrm{\Pi }_{k=1}^n\frac{e^{\frac{i}{2}(\alpha _k+i0)}}{(1e^{i(\alpha _k+i0)})}$$
on $`_\alpha ^n.`$ Since its singular support is the union of the hyperplanes $`_{km}:=\{(\alpha _1,\mathrm{},\alpha _n)^n:\alpha _k=2\pi m\},T`$ has smooth localization to neighborhoods of Floquet exponents $`(\alpha _1,\mathrm{},\alpha _n)`$ which are independent of $`\pi `$ over $`𝐐`$. Hence the regularized trace is simply the evaluation of the distribution trace at a regular point. Similarly, the coefficient of $`D_{s}^{}{}_{}{}^{k1}`$ in (0.12) has a distribution trace of the form
$`(0.14)`$
$$\underset{q𝐍^n}{}_{k,1}(q_1+\frac{1}{2},\mathrm{},q_n+\frac{1}{2})e^{i_{k=1}^n(q_k+\frac{1}{2})\alpha _k}$$
$$=_{k,1}(D_{\alpha _1},\mathrm{},D_{\alpha _n})T(\alpha )$$
for a certain polynomial $`_{k,1}`$. Hence this trace is also a locally smooth function in neighborhoods of non-resonant exponents.
Theorem C The wave invariants are given by
$$a_{k\gamma }=_{k,1}(D_{\alpha _1},\mathrm{},D_{\alpha _n})T(\alpha ).$$
The coefficients of the polynomials $`_{k,1}`$ are evidently polynomials in the QBNF coefficients and the differentiation process produces polynomials in the $`\beta _j`$’s. Combined with Theorem B and a dimensional analysis, this proves Theorem A.
The proofs of Theorems A-C also lead to a somewhat simpler proof of Guillemin’s inverse theorem that the classical (in fact the full quantum) normal form is determined by the wave trace invariants for all the iterates of $`\gamma .`$ We only sketch the proof here, assuming the reader’s familiarity with the original proof of Guillemin in \[G.2\]. The key point is to focus on the Floquet invariants $`\beta _j:=(1e^{\theta _j})^1`$ for all the iterates $`\gamma ^m`$ of $`\gamma `$, that is the residues (0.10) for $`t=L,2L,3L,\mathrm{}.`$ It follows from the calculations in Theorems A-C that the wave invariants are polynomials in the $`\beta _j`$’s: more precisely, for each m, $`a_{k\gamma ^m}`$ is the special value at $`\theta _j=m\alpha _j`$ of the fixed polynomial $`I_{\gamma ;k}`$ in the variables $`\beta _j`$. Under the irrationality condition above, the points $`(e^{im\alpha _1},\mathrm{},e^{im\alpha _n})`$ form a dense set on the torus, and hence the special values at these points determine the entire polynomial. The coefficients of the kth polynomial $`I_{\gamma ;k}`$ are therefore determined by the wave invariants for $`\gamma ,\gamma ^2,\mathrm{}`$. By studying the relation of the coefficients of $`I_{\gamma ;k}`$ to the normal form invariants, Guillemin proves that all of the latter can be determined from the former.
Although Theorems A-C are only proved here under the hypothesis that $`\gamma `$ is non-degenerate elliptic, they have analogues for hyperbolic and mixed hyperbolic-elliptic geodesics, which we plan to describe in a future article \[Z.3\]. We also note that for closed geodesics possessing neighborhoods in which the metric has no pairs of conjugate points for $`tL`$, the wave invariants can be calculated directly from a Hadamard parametrix \[D\] \[Z.1,5\]. Since a sufficiently large number of iterates of an elliptic closed geodesics will always contain pairs of conjugate points , but a small iterate may contain none, the calculation here and in \[Z.1\] overlap but are independent. The calculations of \[Z.1\] also apply to hyperbolic geodesics without pairs of conjugate points, showing that the form of the wave invariants is essentially the same for the hyperbolic and elliptic cases. However the form resulting from the Hadamard parametrix is not immediately that of FJF polynomials, and it takes considerable manipulation to show that the formulae given here and in that paper agree. In the opposite extreme of Zoll manifolds, all of whose geodesics are closed and of completely degenerate elliptic type, the wave invariants are calculated in \[Z3,4\] by yet another method.
The organization of this paper is as follows:
§1: The models
§2: Semi-classical normal form of the Laplacian
§3: Normal form: Proof of Theorem B
§4: Residues and wave invariants: Proof of Theorem C
§5: Local formulae for the residues: Proof of Theorem A
§6: Quantum Birkhoff normal form coefficients
§7: Explicit formulae in dimension 2
§8: Appendix: The classical Birkhoff normal form
§9: Index of Notation
The author wishes to thank S.Graffi for discussions of quantum Birkhoff normal forms during a visit to the University of Bologna in June, 1994, where this work was begun. He also wishes to thank Y.Colin de Verdiere for his remarks on a preliminary version of the results which were presented during a visit to the Institut Fourier in January 1995. The final version was completed after the author received a copy of the article \[G.2\] of Guillemin, and has benefited a good deal from the discussion there of normal forms.
## 1. The models
As mentioned above, the calculation of the wave invariants associated to a closed geodesic $`\gamma `$ of a Riemannian manifold $`(M,g)`$ can be transplanted to the normal bundle $`N_\gamma `$ by means of the exponential map. Thus, the model space is the cylinder $`S_L^1\times ^n`$, where as above $`S_L^1=/L`$. In this section we first collect together some basic formulae and facts concerning the “Hermite package” on this model space: that is, those aspects of analysis which come from the representation theory of the Heisenberg and metaplectic algebras. We then transfer the Hermite package to $`N_\gamma `$ in a way particularly well-adapted to the metric along $`\gamma `$.
§1.1: The model: $`=H^2(S_L^1)L^2(^n)`$
Since we are only concerned with the conic neighborhood (0.2) of $`^+\gamma `$, we only consider the positive part $`T_+^{}S_L^1\times T^{}^n`$ and its quantum analogue the Hardy space $`H^2(S_L^1)L^2(^n)`$ with $`k0.`$
On the phase space level, the model is roughly $`T^{}(S_L^1\times ^n)`$. More precisely it is the cone (0.1) in the natural symplectic coordinates $`(s,\sigma ,y,\eta )`$. Since (0.1) is a conic neighborhood of $`^+\gamma `$, we will view it as a subcone of the positive part $`T_+^{}S_L^1\times T^{}^n`$ ($`\sigma >0`$). On the Hilbert space level the model is then $`:=H^2(S_L^1)L^2(^n)`$, where $`H^2(S_L^1)`$ is the Hardy space, or more precisely the range $`_\psi `$ of the microlocal cutoff $`\psi `$ of the introduction; generally we omit the subscript unless we need to emphasize the role of $`\psi `$. We now introduce some distinguished algebras of operators on the model space.
First is the (complexified) Heiseberg algebra $`𝐡_n`$, which will be identified with its usual realization on $`L^2(^n)`$. It is then generated by the elements $`y_j=`$ “multiplication by $`y_j`$” and by $`D_{y_j}=\frac{}{iy_j}`$, or equivalently by the annihilation, resp. creation, operators
$$A_j:=y_j+iD_{y_j}A_j^{}=y_jiD_{y_j}$$
which satisfy the commutation relations
$$[A_j,A_k]=[A_j^{},A_k^{}]=0[A_j,A_k^{}]=2\delta _{ij}I.$$
The enveloping algebra of the Heisenberg algebra
$`(\mathrm{1.1.1})`$
$$:=<Y_1,\mathrm{},Y_n,D_1,\mathrm{},D_n>$$
is the algebra of partial differential operators on $`^n`$ with polynomial coefficients. We let $`^n`$ denote the subspace of polynomials of degree n in the variables $`y_j,D_{y_j}.`$ In the usual isotropic Weyl algebra $`𝒲^{}`$ of pseudo-differential operators on $`^n`$, the operators $`y_j,D_{y_j}`$ are given the order $`\frac{1}{2}`$, so that
$`(\mathrm{1.1.1}b)`$
$$^n𝒲^{n/2}$$
$$[^m,^n]^{m+n2}.$$
The symplectic algebra $`sp(n,)`$ is represented in $`^2`$ by homogeneous quadratic polynomials in $`Y_j,D_j`$, and a maximal abelian subalgebra of it is spanned by the harmonic oscillators (0.2). We denote by
$`(\mathrm{1.1.2}a)`$
$$:=<I_1,\mathrm{},I_n>$$
the (maximal abelian) subalgebra they generate in $`𝒲`$, with $`^k:=𝒲^k`$, and by
$`(\mathrm{1.1.2}b)`$
$$𝒫_{}=$$
the subalgebra of polynomials in the generators (0.2)), with $`𝒫_{}^k`$ the space of polynomials of degree k.
The full pseudo-differential algebra on $`S_L^1\times ^n`$ is the doubly filtered algebra
$$\mathrm{\Psi }^{}(S_L^1\times ^n)\mathrm{\Psi }^{}(S_L^1)𝒲^{}$$
with
$$\mathrm{\Psi }^{mn}(S_L^1\times ^n)\mathrm{\Psi }^m(S_L^1)𝒲^n.$$
A maximal abelian subalgebra of it is given by
$`(\mathrm{1.1.3})`$
$$𝒜:=<D_s,I_1,\mathrm{},I_n>=<,I_1,\mathrm{},I_n>$$
where $``$ is the distinguished element (0.3). It inherits a double filtration $`𝒜^{mn}`$. As above, our interest is really in the microlocalization of (1.1.3) to the cone (0.2), i.e. the operators in (1.1.3) will only be used in composition with the microlocal cut-off $`\overline{\psi }(,I_1,\mathrm{},I_n)`$. In this cone $`D_s`$ and $``$ are elliptic. Hence the subalgebra
$`(\mathrm{1.1.4})`$
$$<>𝒫_{}$$
of pseudo-differential symbols in $``$ with coefficients in $`𝒫_{}`$ is well defined.
An orthonormal basis of $`L^2(^n)`$ of joint eigenfunctions of $``$ is provided by the Hermite functions $`\gamma _q,q^n.`$ Here, $`\gamma _o`$ is the Gaussian $`\gamma _o(y)=\gamma _{iI}(y):=e^{\frac{1}{2}|y|^2}`$. It is the unique “vacuum state”, i.e. the state annihilated by the annihilation operators. The qth Hermite function is then given by $`\gamma _q:=C_qA_1^{q_1}\mathrm{}A_n^{q_n}\gamma _o(q𝐍^n\}),`$ with $`C_q=(2\pi )^{n/2}(q!)^{1/2},q!=q_1!\mathrm{}q_n!.`$ The notation “$`\gamma _q`$” for “Gaussian” is standard in this context, see \[F\], and should not be confused with the notation for closed geodesics. An orthonormal basis of $`H^2(S_L^1)L^2(^n)`$ of joint $`𝒜`$\- eigenfunctions is then furnished by
$`(\mathrm{1.1.5})`$
$$\varphi _{kq}^o(s,y):=e_k(s)\gamma _q(y),e_k(s):=e^{\frac{2\pi }{L}iks}.$$
§1.2: The twisted model $`_\alpha `$
We now introduce a unitarily equivalent (twisted) version of the model, in which the distinguished element $``$ gets conjugated to $`D_s.`$ This will eventually help to simplify the transport equations in §2.
The unitary equivalence will be given by conjugation with the unitary operator
$`(\mathrm{1.2.1}a)`$
$$\mu (r_\alpha ):=_{S_L^1}^{}\mu (r_\alpha (s))𝑑s=_{S_L^1}^{}e^{i\frac{s}{L}H_\alpha }𝑑s$$
where $`\mu `$ is the metaplectic representation, and where $`r_\alpha (s)`$ is the block diagonal orthogonal transformation on $`^{2n}`$ with blocks
$`(\mathrm{1.2.1}b)`$
$$r_{\alpha _j}(s):=\left(\begin{array}{cc}cos\alpha _j\frac{s}{L}\hfill & sin\alpha _j\frac{s}{L}\hfill \\ sin\alpha _j\frac{s}{L}\hfill & cos\alpha _j\frac{s}{L}\hfill \end{array}\right).$$
The direct integral here refers to the representation of $`L^2(S^1)L^2(^n)`$ as $`_{S_L^1}L^2(^n)𝑑s`$, that is,
$$_{S_L^1}^{}\mu (r_\alpha (s))𝑑sf(s,y)=\mu (r_\alpha (s))f(s,y)$$
where the right side is the application the operator in the $`y`$-variables with $`s`$ fixed. As will be seen below, $`\mu (r_\alpha )`$ conjugates $``$ to $`D_s`$, and commutes with $`I_1,\mathrm{},I_n`$. Hence it preserves the algebra (1.1.3).
On the other hand, it does not preserve the Hilbert space $`.`$ Indeed, the elements $`\varphi _{kq}^o`$ get transformed into the elements
$`(\mathrm{1.2.2})`$
$$e_{kq}(s)\gamma _q(y),e_{kq}(s):=e^{ir_{kq}s}$$
which are not periodic in $`s`$. Rather they satisfy
$`(\mathrm{1.2.3}a)`$
$$e_{kq}\gamma (s+L,y)=e^{i\kappa _q}e_{kq}\gamma (s,y)$$
with
$`(\mathrm{1.2.3}b)`$
$$\kappa _q=\underset{j=1}{\overset{n}{}}(q_j+\frac{1}{2})\alpha _j.$$
The space $`H^2(S_L^1)L^2(^n)`$ thus gets taken to the space $`_\alpha `$ of elements of the form
$`(\mathrm{1.2.4})`$
$$f(s,y)=\underset{k=o}{\overset{\mathrm{}}{}}\underset{q^n}{}\widehat{f}(k,q)e_{kq}\gamma _q$$
with square summable coefficients.
We can better describe this Hilbert space (and its associated phase space) in the language of ‘quantized mapping cylinders’.
On the phase space level we have the symplectic map $`r_\alpha (L)`$ of $`T^{}^n`$ of (1.2.1b), essentially the Poincare map of our problem. As in \[F.G\]\[G.2\] we can introduce its (homogeneous) symplectic mapping cylinder $`C_{r_\alpha (L)}`$: namely, the quotient of $`T^{}\times T^{}^n`$ under the cylic group $`<R_\alpha (L)>`$ generated by the symplectic map
$`(\mathrm{1.2.5})`$
$$\stackrel{~}{r}_\alpha (L):T^{}\times T^{}^nT^{}\times T^{}^n,\stackrel{~}{r}_\alpha (L)(s,\sigma ,y,\eta ):=(s+L,\sigma ,r_\alpha (L)(y,\eta )).$$
Note that the first return time is constant, which is consistent with \[F.G\]\[G.2, 2.10\] since elements of $`Sp(2n,)`$ preserve the contact form $`ds\frac{1}{2}(\eta dyyd\eta ).`$ Also note that the mapping cylinder can be untwisted via the symplectic map
$`(\mathrm{1.2.7})`$
$$R_\alpha :T^{}\times T^{}^nT^{}\times T^{}^n,R_\alpha (s,\sigma ,y,\eta ):=(s,\sigma +\frac{1}{2}\underset{j}{}\alpha _j(y_j^2+\eta _j^2),r_\alpha (s)(y,\eta )).$$
Indeed, we have
$`(\mathrm{1.2.8})`$
$$R_\alpha (s+L,\sigma ,y,\eta )=\stackrel{~}{r}_\alpha (L)R_\alpha (s,\sigma ,y,\eta )$$
so that $`R_\alpha `$ induces a symplectic equivalence $`C_{r_\alpha (L)}T^{}(S_L^1)\times R^{}(^n).`$
On the quantum level, the analogue of the mapping cylinder is the Hilbert space $`_\alpha `$ of functions on $`\times ^n`$ satisfying
$`(\mathrm{1.2.6})`$
$$f(s+L,y)=\mu (r_\alpha (L))f(s,y)$$
and square integrable on $`[0,L)\times ^n`$. The intertwining operator $`\mu (r_\alpha )`$ is essentially the analogue of $`R_\alpha `$. More precisely, we have:
(1.2.9) Proposition
(i) $`\mu (r_\alpha )^{}D_s\mu (r_\alpha )=;`$
(ii) $`\mu (r_\alpha )\varphi _{kq}^o(s,y)=e^{ir_{kq}s}\gamma _q;`$
(iii) $`\mu (r_\alpha ):_\alpha `$.
Proof:
(i) Follows from the fact that $`\mu `$ takes the jth diagonal block
$$\left(\begin{array}{cc}0\hfill & 1\hfill \\ 1\hfill & 0\hfill \end{array}\right)$$
to $`I_j`$ and hence that $`\mu (r_\alpha (s))=e^{i\frac{s}{L}H_\alpha }`$.
(ii) Follows from (i) and the fact that $`\gamma _q`$ is an eigenfunction of eigenvalue $`q_j+\frac{1}{2}`$ of $`I_j`$.
(iii) Follows from (ii).∎
Under conjugation by $`\mu (r_\alpha )`$, the sub-algebra (1.1.4) goes over to the algebra
$`(\mathrm{1.2.10})`$
$$<D_s>𝒫_{}$$
of pseudodifferential symbols in $`D_s`$ with coefficients in the $`I_j`$’s. As usual, it is understood to be microlocalized to the cone (0.2).
We now conjugate by a further unitary equivalence to transfer this Hermite package to $`N_\gamma `$. It will induce a new model better adapted to the geometry of $`(M,g)`$ near $`\gamma ;`$ we call it the adapted model.
§1.3: The adapted model
To define it, we must discuss the Jacobi equation along $`\gamma .`$ Let Y(s) be a vector field along $`\gamma ,`$ and as above let $`Y(s)=_{i=1}^nY_i(s)e_i(s)`$ be its expression in terms of the parallel normal frame. The Jacobi equation is then
$$\frac{d^2}{ds^2}Y_i+\underset{ij=1}{\overset{n}{}}K_{ij}Y_j=0.$$
Let $`𝒥_\gamma ^{}𝐂`$ denote the space of complex orthogonal Jacobi fields along $`\gamma .`$ Equipped with the symplectic form $`\omega (X,Y):=<X,\frac{DY}{ds}><\frac{DX}{ds},Y>,`$ it is a symplectic vector space of dimension 2n (over $`𝐂`$). Here $`\frac{D}{ds}`$ denotes covariant differentiation along $`\gamma ,`$ and $`<,>`$ is the inner product defined by the metric. Now let $`Z(s):=(Y(s),\frac{DY}{ds}).`$ The linear Poincare map $`P_\gamma `$ is by defintion the operator on $`𝒥_\gamma ^{}𝐂`$ given by $`P_\gamma Z(t)=Z(t+L).`$ We recall that $`\gamma `$ is assumed non-degenerate elliptic, hence
$$Spec(P_\gamma )S^1\pm 1spec(P_\gamma ).$$
Since $`P_\gamma Sp(2n,𝐑),`$ its eigenvalues come in complex conjugate pairs $`\{e^{i\alpha _k},e^{i\alpha _k}\}.`$ The exponents $`\alpha _k`$ (called Floquet) are only defined up multiples of $`2\pi ;`$ they will be normalized below as in \[B.B, (9.3.17)\]. Then there exists a basis of complex eigenvectors $`\{Y_1,\mathrm{},Y_n\}`$ satisfying
$`(\mathrm{1.3.1})`$
$$P_\gamma Y_k=e^{i\alpha _k}Y_kY_j𝒥𝐂$$
$$\omega (Y_k,Y_j)=\omega (\overline{Y}_j,\overline{Y}_k)=0,\omega (Y_j,\overline{Y}_k)=\delta _{jk}.$$
with one choice of eigenvalue from each complex conjugate pair. Equivalently, the span of $`\{Y_1,\mathrm{},Y_n\}`$ defines a $`P_\gamma `$-invariant positive Lagrangean subspace of $`𝒥_\gamma ^{}𝐂`$.
Consider now the (modified) Wronskian matrix
$`(\mathrm{1.3.2})`$
$$a_s:=\left(\begin{array}{cc}Im\dot{Y}(s)^{}\hfill & ImY(s)^{}\hfill \\ Re\dot{Y}(s)^{}\hfill & ReY(s)^{}\hfill \end{array}\right).$$
Here, with a little abuse of notation, we denote the components of $`Y(s)`$, resp. $`\frac{DY}{ds}`$ relative to the parallel normal frame $`e(s)`$ by $`Y(s)`$, resp. $`\dot{Y}(s)`$ and the notation $`A^{}`$ refers to the adjoint of a matrix $`A`$. From (1.3.1) we have that $`a_sSp(2n,𝐑)`$ , which is equivalent to
$$\overline{Y}\frac{DY}{ds}\overline{\frac{DY}{ds}}Y=iI.$$
As above, we let $`\mu `$ denote the metaplectic representation. Identifying $`a_s`$ with one of its two possible lifts from $`Sp(2n,)`$ to $`Mp(2n,R)`$ we introduce the unitary operator
$`(\mathrm{1.3.3})`$
$$\mu (a):=_\gamma ^{}\mu (a_s)𝑑s$$
on $`_{S_L^1}^{}L^2(^n)𝑑s`$. In other words,
$$\mu (a)f(s,y)=\mu (a_s)f(s,y)$$
where the operator on the right side acts in the $`y`$-variables. Informally, we think of the range as the Hilbert space $`_\gamma ^{}L^2(N_{\gamma (s)})𝑑s,`$ and below we will describe it more completely in terms of quantum mapping cylinders. We now use $`\mu (a)`$ to transfer the Hermite package to $`N_\gamma `$. We begin with the generators of the Heisenberg algebra, and set:
$$\begin{array}{cc}P_j:=\mu (a)^{}D_j\mu (a)\hfill & Q_j:=\mu (a)^{}Y_j\mu (a)\hfill \\ \mathrm{\Lambda }_j:=\mu (a)^{}A_j\mu (a)\hfill & \mathrm{\Lambda }_j^{}=\mu (a)^{}A_j^{}\mu (a)\hfill \end{array}.$$
$$I_{\gamma j}:=\frac{1}{2}\mathrm{\Lambda }_j\mathrm{\Lambda }_j^{}.$$
We will refer to the $`\mathrm{\Lambda }_j`$’s, resp. $`\mathrm{\Lambda }_j^{}`$’s, as the adapted annihilation and creation operators and to the operators $`I_{\gamma j}`$ in the symplectic algebra as the adapted action operators. These adapted operators play a key role in the study of quasi-modes associated to $`\gamma `$. To establish the connection, we now verify that they coincide with the operators similarly denoted in \[B.B., ch.9\].
First, some remarks on notation. For reasons that will become clearer below (§1.4-5), we change the notation for the transverse coordinates from $`y`$ to $`u`$, which should be thought of as the rescaled coordinates $`u=L^{\frac{1}{2}}y`$. Objects in the adapted model will henceforth always be expressed in terms of $`u`$coordinates. For instance, multiplication by $`u_j`$ will be denoted simply by $`u_j`$ and differentiation in $`u_j`$ by $`D_{u_j}`$. Also, to quote easily some basic facts about the metaplectic representation from Folland \[F\], we will conform to the following ‘transposed’ notation for the remainder of this section: symbols of operators will be denoted $`\sigma (p,q)`$ rather than $`\sigma (q,p)`$, and the corresponding Weyl pseudodifferential operator will be denoted by $`\sigma (D,x)`$ (later we will also use the more standard notation $`\sigma ^w(x,D)`$). An element of $`Sp(n,)`$ will be denoted by $`𝒜`$ and its transpose by $`𝒜^{}`$.
(1.3.4) Proposition
$$\begin{array}{cc}\mathrm{\Lambda }_j=_{k=1}^n(iy_{jk}D_{u_k}\frac{dy_{jk}}{ds}u_k)\hfill & \mathrm{\Lambda }_j^{}=_{k=1}^n(i\overline{y}_{jk}D_{u_k}\overline{\frac{dy_{jk}}{ds}}u_k)\hfill \end{array}$$
Proof: This follows from the metaplectic covariance of the Weyl calculus, that is from the identity \[F, Theorem 2.15\]
$$(\sigma 𝒜)(D,x)=\mu (𝒜^{})\sigma (D,x)\mu (𝒜^{})^1.$$
If we set $`a=𝒜^{}`$ we get (in an obvious vector notation)
$$\mu (a)^{}(u_j+iD_{u_j})\mu (a)=i(Y_jD_u)\dot{Y}_ju$$
(compare \[BB, ch.9\]). ∎
As for the tangential operator, we have:
(1.3.5) Proposition The image $``$ of $`D_s`$ under $`\mu `$ is given by:
$$:=\mu (a)^{}D_s\mu (a)=D_s\frac{1}{2}(\underset{j=1}{\overset{n}{}}D_{u_j}^2+\underset{ij=1}{\overset{n}{}}K_{ij}(s)u_iu_j).$$
Proof: The left side is equal to $`(D_s+\mu (a_s)^{}D_s\mu (a_s)).`$ To evaluate the second term, we use that both ReY(s) and ImY(s) are Jacobi fields, and that Jacobi’s equation is equivalent to the linear system $`\frac{D}{ds}(Y,P)=JH(Y,P).`$ Here, $`P=\frac{DY}{ds},`$ J is the standard complex structure on $`R^{2n},`$ and
$$H=\left(\begin{array}{cc}K\hfill & 0\hfill \\ 0\hfill & I\hfill \end{array}\right)$$
where $`K`$ is the curvature matrix and $`I`$ is the identity matrix \[B.B, (9.2.9)\]. Hence, the second term is $`\frac{1}{i}d\mu (JH)`$ with $`d\mu `$ the derived metaplectic representation. But $`\frac{1}{i}d\mu (JH)=1/2(_{i=1}^n_{u_i}^2_{ij=1}^nK_{ij}(s)u_iu_j)`$ \[F\]. ∎
We now consider the appropriate Hilbert space (quantized mapping cylinder) in the adapted model. We first note:
(1.3.6) Proposition
(i)
$$\mu (a^1)\gamma _o(s,u):=U_o(s,u)=(detY(s))^{1/2}e^{i\frac{1}{2}\mathrm{\Gamma }(s)u,u}$$
where $`\mathrm{\Gamma }(s):=\frac{dY}{ds}Y^1.`$
(ii)
$$\mu (a^1)\gamma _q:=U_q=\mathrm{\Lambda }_1^{q_1}\mathrm{}\mathrm{\Lambda }_n^{q_n}U_o.$$
Proof: Let us recall the action of the metaplectic representation on Gaussians \[F, Ch. 4.5\]. For $`Z`$ in the Siegel upper half space (n x n complex matrices $`Z=X+iY`$ with $`Y>>0`$), we have the Gaussian
$$\gamma _Z(x):=e^{i\frac{1}{2}<Zx,x>}$$
on $`^n.`$ The action of an element $`𝒜Mp(2n,)`$ on the Gaussian is given by
$$\mu (𝒜^1)\gamma _Z=m(𝒜,Z)\gamma _{\alpha (𝒜)Z}$$
where
$$m(𝒜,Z)=det^{1/2}(CZ+D),\alpha (𝒜)Z=(AZ+B)(CZ+D)^1$$
for
$$𝒜=\left(\begin{array}{cc}A\hfill & B\hfill \\ C\hfill & D\hfill \end{array}\right)$$
(see \[F, 4.65\].) Writing
$$Y(s)=ImY(s)i+ReY(s),\mathrm{\Gamma }(s)=(Im\dot{Y}(s)i+Re\dot{Y})(ImY(s)i+ReY(s))^1$$
we see that $`m(𝒜,iI)\gamma _{\alpha (𝒜)iI}(u)=(detY(s))^{1/2}exp(\frac{i}{2}<\mathrm{\Gamma }(s)u,u>`$ if
$$𝒜=\left(\begin{array}{cc}Im\dot{Y}(s)\hfill & Re\dot{Y}\hfill \\ ImY(s)\hfill & ReY(s)\hfill \end{array}\right).$$
The formula (i) follows from $`𝒜=a^{}`$ and (ii) is an immediate consequence of (i).∎
Consider now the periodicity properties of the above data under $`ss+L`$. We observe that $`a_s`$ fails to be periodic in $`s`$ for two reasons: first, due to the holonomy of the frame $`e(s)`$, and secondly due to the monodromy of the Jacobi fields $`Y(s)`$. Indeed, we have:
$`(\mathrm{1.3.7}a)`$
$$a_{(s+L)}=Ta_sr_\alpha ^1(L)$$
with $`r_\alpha `$ as in §1.2 and where $`T`$, the holonomy matrix, is the 2n by 2n block diagonal matrix with equal diagonal blocks $`t=(t_{ij})`$ satisfying
$$e_i(L)=\underset{j=1}{\overset{n}{}}t_{ij}e_j(0).$$
It is of course the lift to $`T^{}^n`$ of a rotation on the base. The two properties can be summarized by writing
$`(\mathrm{1.3.7}b)`$
$$\stackrel{~}{a}_{s+L}=T\stackrel{~}{a}_s$$
with $`\stackrel{~}{a}_s:=a_sr_\alpha (s)`$.
As in §1.2 we will reformulate (1.3.7a-b) in terms of quantum mapping cylinders. First, we put
$`(\mathrm{1.3.8})`$
$$C_T^{\mathrm{}}(\times ^n):=\{fC^{\mathrm{}}(\times ^n):f(s+L,u)=\mu (T)f(s,u)\}$$
and let $`_T`$ denote its closure with respect to the obvious inner product over $`[0,L)\times ^n`$. Note that the metaplectic operator $`\mu (T)`$ is simply
$$\mu (T)f(u)=f(t^1u)$$
and hence that
$$C_T^{\mathrm{}}(\times ^n)C^{\mathrm{}}(N_\gamma )$$
where the isomorphism is simply the pull-back by the exponential map defined by the frame $`e(s)`$. Thus:
(1.3.8)Proposition
(i) Let $`P(s,D_s,u,D_u)`$ be a partial differential operator on $`N_\gamma `$, expressed in the coordinates $`(s,u).`$ Then
$$P(s+L,D_s,u,D_u)=\mu (T)P(s,D_s,u,D_u)\mu (T)^{};$$
(ii) The functions
$$\mu (\stackrel{~}{a}_s)(\varphi _{kq}^o)=\varphi _{kq}:=e^{ir_{kq}s}U_q(s,u)$$
define a smooth orthonormal basis of $`_T`$.
Proof
(i) It suffices to prove this when $`P`$ is a vector field given in the local normal coordinates by $`a_o(s,u)D_s+_{j=1}^na_j(s,u)D_{u_j}.`$ Since the metaplectic operator $`\mu (T)`$ corresponding to $`T`$ is the operator $`f(u)f(t^1u),`$ we have
$$\mu (T)P\mu (T))^{}=a_o(s,t^1u)D_s+_{j=1}^na_j(s,t^1u)t_{jk}D_{u_k}.$$
On the other hand, the vector field is well-defined on $`N_\gamma `$ if and only if
$$a_o(s+L,t^1u)D_s+a_i(s+L,t^1u)t_{ij}D_{u_j}=a_o(s,u)D_s+a_j(s,u)D_{u_j}.$$
(ii) Clear, since by (1.3.7b) $`\mu (\stackrel{~}{a})`$ intertwines the model and the quantum mapping cylinder of $`\mu (T).`$
Remark: Statement (ii) is equivalent to
$$U_q(s+L,t^1u)=e^{i\kappa _q}U_q(s,u)$$
(correcting the formula stated in \[B.B., (9.3.25)\].)
It follows that we may write a smooth function $`f`$ in the adapted model in the form
$`(\mathrm{1.3.9})`$
$$f(s,u)=\underset{k=o}{\overset{\mathrm{}}{}}\underset{q𝐍^n}{}\widehat{f}(k,q)e^{ir_{kq}s}U_q(s,u).$$
§1.4 Metric scaling and weights
As mentioned above, $`\mathrm{\Delta }`$ and hence the wave invariants have well-defined weights under the metric rescaling $`gϵ^2g`$. Since the wave invariants will be expressed in terms of QBNF coefficients, it is natural to ask how the latter scale. The question is not really well-posed since the QBNF coefficients are coefficients with repsect to Harmonic oscillators $`_{y_j}^2+y_j^2`$ whose scaling behaviour depends on the choice of coordinates. To amplify this point, we record how various metric objects scale under metric re-scaling.
In the following table, $`(s,y)`$ denote the Fermi normal coordinates relative to $`g`$, $`(s,u)`$ denote the scaled coordinates $`(s,L^{\frac{1}{2}}y)`$ and $`p_u`$ denotes the symplectic coordinates dual to $`u`$.
| $`g`$ | $`ϵ^2g`$ |
| --- | --- |
| $`L,(s,y)`$ | $`ϵL,(ϵs,ϵy)`$ |
| $`_s,_{y_j},e_j`$ | $`ϵ^1_s,ϵ^1_{y_j},ϵ^1e_j`$ |
| $`y_{ij}:=g(Y_i,e_j)`$ | $`ϵ^{\frac{1}{2}}y_{ij}=ϵ^2g(ϵ^{\frac{1}{2}}Y_j,ϵ^1e_j)`$ |
| $`K_{ij}=g(R(_s,e_i)_s,e_j)`$ | $`ϵ^2K_{ij}`$ |
| $`u=L^{\frac{1}{2}}yp_u=L^{\frac{1}{2}}\eta `$ | $`ϵ^{\frac{1}{2}}uϵ^{\frac{1}{2}}p_u`$ |
| $`\mathrm{\Lambda }_j=_{k=1}^n(iy_{jk}D_{u_k}\dot{y}_{jk}u_k)`$ | $`_{k=1}^n(iy_{jk}D_{u_k}\dot{y}_{jk}u_k)`$ |
| $`:=D_s\frac{1}{2}(_{j=1}^nD_{u_j}^2+_{ij=1}^nK_{ij}(s)u_iu_j)`$ | $`ϵ^1[D_s\frac{1}{2}(_{j=1}^nD_{u_j}^2+_{ij=1}^nK_{ij}(s)u_iu_j)]`$ |
The entry $`y_{ij}ϵ^{\frac{1}{2}}y_{ij}`$ follows from the scale invariance of the Jacobi equation together with the normalization condition
$$g(ReY_i,Im\dot{Y}_i)g(Re\dot{Y}_i,ImY_i)=Constant.$$
which implies that $`Y_iϵ^{\frac{1}{2}}Y_i`$.
We observe that the creation/annihilation operators, hence the harmonic oscillators $`\mathrm{\Lambda }_j^{}\mathrm{\Lambda }_j`$, of the adapted model are scale-invariant, and that the distinguished element $``$ has weight -1. These are the desired scaling properties and we would like the basic and twisted models to possess them as well. As they stand, these models do not scale properly if we interpret the $`(s,y)`$-coordinates as Fermi normal coordinates. However, they do scale properly if we interpret the $`y`$ coordinates as weightless.
§1.5 Scaled adapted model and intertwining operators
To avoid confusion, we now introduce the weightless coordinates $`x=L^1y=L^{\frac{1}{2}}u,\xi =L\eta =L^{\frac{1}{2}}p_u`$ and henceforth use them exclusively for the scaled adapted, basic and twisted models. In the following table we record how the various objects appear in the weightless coordinates. We also record the various intertwining operators, since they get altered when we used weightless coordinates. For instance, intertwining by $`\mu (r_\alpha )`$ above is weightless but that by $`\mu (a)`$ is of mixed weight.
| Adapted model | Scaled adapted model |
| --- | --- |
| $`\mathrm{\Lambda }_j`$ | $`_{k=1}^n(iy_{jk}L^{\frac{1}{2}}D_{x_k}L^{\frac{1}{2}}\frac{dy_{jk}}{ds}x_k)`$ |
| $``$ | $`D_s\frac{1}{2}(_{j=1}^nL^1D_{u_j}^2+_{ij=1}^nLK_{ij}u_iu_j)`$ |
| $`U_o(s,u)`$ | $`U_{Lo}(s,x)=(detY(s))^{1/2}e^{\frac{i}{L}\mathrm{\Gamma }(s)x,x}|dx|^{\frac{1}{2}}`$ |
The following intertwining operators will arise in the construction of the normal form:
| From * to * | Classical | Quantum |
| --- | --- | --- |
| Ad.model to Sc.ad.model | $`(u,p_u)(L^{\frac{1}{2}}u,L^{\frac{1}{2}}p_u)=(x,\xi )`$ | $`\mu (D_{L^{\frac{1}{2}}})`$ |
| Sc.Ad.Model to Tw.model | $`(x,\xi )(L^{\frac{1}{2}}ReYx+L^{\frac{1}{2}}ImY\xi ,L^{\frac{1}{2}}Re\dot{Y}x+L^{\frac{1}{2}}Im\dot{Y}\xi )`$ | $`\mu (D_{L^{\frac{1}{2}}}a)`$ |
Above, the notation $`\mu (D_r)`$ refers to the metaplectic (dilation) operator corresponding to the symplectic matrix
$$D_r:=\left(\begin{array}{cc}r\hfill & 0\hfill \\ 0\hfill & r^1\hfill \end{array}\right),\mu (D_r)f(s,y):=f(s,r^1y).$$
## 2. Semi-classical normal form of the Laplacian
We return now to $`\sqrt{\mathrm{\Delta }}`$, which as in the introduction will be identified with its transfer to $`N_\gamma `$ under the exponential map. The finite jets of this transfer are globally well- defined on $`N_\gamma `$, so we will often treat $`\mathrm{\Delta }`$ as if it too were globally well-defined. Our purpose is to define the semi-classically rescaled Laplacian $`\mathrm{\Delta }_h`$ and to put $`\mathrm{\Delta }_h`$ into a semi-classical normal form. This is the crucial preliminary step in putting $`\mathrm{\Delta }`$ itself into normal form.
In view of §1.5, there are two rescalings at hand: the semi-classical rescaling $`\mathrm{\Delta }_h`$ and the metric $`L`$-rescaling above. The two rescalings have quite distinct origins, so we have kept them separate.
To motivate the rescalings and the emergence of semi-classical asymptotics, let us recall that the quasi-modes associated to $`\gamma `$ have the form
$`(2.1)`$
$$\mathrm{\Phi }_{kq}(s,\sqrt{r_{kq}}y)=e^{ir_{kq}s}\underset{j=0}{\overset{\mathrm{}}{}}r_{kq}^{\frac{j}{2}}U_q^{\frac{j}{2}}(s,\sqrt{r_{kq}}y,r_{kq}^1)$$
with $`U_q^o=U_q`$ (see \[B.B\]). The intertwining operator $`W_\gamma `$ to the normal form is then the operator defined by the equations
$`(2.2)`$
$$W_\gamma \varphi _{kq}(s,y)=\mathrm{\Phi }_{kq}(s,\sqrt{r_{kq}}y).$$
The higher order terms $`U^{\frac{j}{2}}`$, hence $`W_\gamma `$, are determined by the conditions
$`(2.3)`$
$$\mathrm{\Delta }_ye^{ir_{kq}s}U_{kq}(s,\sqrt{r_{kq}}y,r_{kq}^1)\lambda _{kq}e^{ir_{kq}s}U_{kq}(s,\sqrt{r_{kq}}y,r_{kq}^1)$$
with $`\lambda _{kq}`$ given by (0.9).
Now write
$$r_{kq}=\frac{1}{h_{kq}L},h_{kq}:=(2\pi k+\underset{j=1}{\overset{n}{}}(q_j+\frac{1}{2}\alpha _j))^1$$
so that the Planck constants $`h_{kq}`$ are metric-independent. In the scaled Fermi coordinates $`u`$ of §1.3-5, the quasi-modes have then the form
$`(2.1^{})`$
$$\mathrm{\Phi }_{kq}(s,h^{\frac{1}{2}}u)=e^{i\frac{s}{h_{kq}L}}U_{kq}(s,\sqrt{h_{kq}}^{\frac{1}{2}}u,h_{kq})$$
and the eigenvalue problem (2.3) becomes
$`(2.3^{})`$
$$\mathrm{\Delta }_ue^{\frac{i}{h_{kq}L}s}U(s,h_{kq}^{\frac{1}{2}}u,h_{kq})=\lambda (h_{kq})e^{\frac{i}{h_{kq}L}s}U(s,h_{kq}^{\frac{1}{2}}u,h_{kq}).$$
We are thus led to study the asymptotic eigenvalue problem
$`(2.4)`$
$$\mathrm{\Delta }_ue^{\frac{i}{hL}s}U(s,h^{\frac{1}{2}}u,h)=\lambda (h)e^{\frac{i}{hL}s}U(s,h^{\frac{1}{2}}u,h)$$
on $`C^{\mathrm{}}(^1\times ^n)`$ with $`U(s,u,h)`$ and $`\lambda (h)`$ asympotic series in $`h`$. Since $`\mathrm{\Delta }_u`$ comes from an operator on $`N_\gamma `$, the eigenvalue problem is taking place on the quantized mapping cylinder $`_T`$ (§1.3) of the adapted model.
We also note that the local expression for the Laplacian in the $`(s,u)`$ coordinates is the same as in the $`(s,y)`$ (Fermi) coordinates:
$$\mathrm{\Delta }_u=\frac{1}{\sqrt{g}}\underset{ij=o}{\overset{n}{}}_{u_i}g^{ij}\sqrt{g}_{u_j}=\mathrm{\Delta }_y|_{u_j:=L^{\frac{1}{2}}y_j}$$
since $`g^{ij}L^1g^{ij}`$ and $`_{u_i}L^{\frac{1}{2}}_{y_i}.`$ Momentarily we are going to rescale the coordinates again to the weightless $`x`$-coordinates of §1.5, and again the Laplacian will be given by the usual expression. Hence it will be a simple matter to pass back and forth between the $`(s,y),(s,u)`$ and $`(s,x)`$ expressions.
It is natural at this point to introduce the unitary operators $`T_h`$ and $`M_h`$ on $`_T`$ or equivalently on the 1/2-density version $`L_T^2(^1\times ^n,\mathrm{\Omega }_{1/2})`$ given by
$`(2.5a)`$
$$T_h(f(s,u)|ds|^{1/2}|du|^{1/2}):=h^{n/2}f(s,h^{\frac{1}{2}}u)|ds|^{1/2}|du|^{1/2}$$
$`(2.5b).`$
$$M_h(f(s,u)|ds|^{1/2}|du|^{1/2}):=e^{\frac{i}{hL}s}f(s,y)|ds|^{1/2}|du|^{1/2}$$
We easily see that:
$`(2.6)`$
$$T_h^{}D_{u_j}T_h=h^{\frac{1}{2}}D_{u_j}$$
$$T_h^{}u_iT_h=h^{\frac{1}{2}}u_i$$
$$M_h^{}D_sM_h=((hL)^1+D_s)$$
$$[M_h,u_i]=[M_h,D_{u_i}]=[M_h,T_h]=[T_h,D_s]=0.$$
Definition The rescaling of an operator $`A_u=a(s,D_s,u,D_u)`$ of the adapted model is the operator
$`(2.7)`$
$$A_h:=T_h^{}M_h^{}AT_hM_h$$
We observe that the operation of rescaling is weightless. In particular, the rescaled Laplacian $`\mathrm{\Delta }_h`$ in the sense of (2.7) is of weight -2. To calculate it, we first note that the (1/2-density) Laplacian in scaled Fermi normal coordinates is given by the expression
$`(2.8)`$
$$\mathrm{\Delta }_u=(J^1_sJg^{oo}_s+\underset{ij=1}{\overset{n}{}}J^1_{u_i}g^{ij}J_{u_j})$$
where $`_x:=\frac{}{x},`$ and where $`J=J(s,u)=\sqrt{g}`$ is the volume density in these coordinates. To obtain a self-adjoint operator with respect to the Lesbegue density $`|ds||du|`$, we replace $`\mathrm{\Delta }`$ by the unitarily equivalent 1/2-density Laplacian
$$\mathrm{\Delta }_{1/2}:=J^{1/2}\mathrm{\Delta }J^{1/2},$$
which can be written in the form:
$`(2.9)`$
$$\mathrm{\Delta }_{1/2}=J^{1/2}_sg^{oo}J_sJ^{1/2}+\underset{ij=1}{\overset{n}{}}J^{1/2}_{u_i}g^{ij}J_{u_j}J^{1/2}$$
$$g^{oo}_s^2+\mathrm{\Gamma }^o_s+\underset{ij=1}{\overset{n}{}}g^{ij}_{u_i}_{u_j}+\underset{i=1}{\overset{n}{}}\mathrm{\Gamma }^i_{u_i}+\sigma _o.$$
From now on, we will only use $`\mathrm{\Delta }_{\frac{1}{2}}`$ and denote it simply by $`\mathrm{\Delta }.`$
We then have:
$`(2.10)`$
$$M_h^{}\mathrm{\Delta }M_h=(hL)^2g^{oo}+2i(hL)^1g^{oo}_s+i(hL)^1\mathrm{\Gamma }^o+\mathrm{\Delta }$$
Conjugation with $`T_h`$ then gives
$`(2.11)`$
$$\mathrm{\Delta }_h=(hL)^2g_{[h]}^{oo}+2i(hL)^1g_{[h]}^{oo}_s+i(hL)^1\mathrm{\Gamma }_{[h]}^o+h^1(\underset{ij=1}{\overset{n}{}}g_{[h]}^{ij}_{u_i}_{u_j})+h^{\frac{1}{2}}(\underset{i=1}{\overset{n}{}}\mathrm{\Gamma }_{[h]}^i_{u_i})+(\sigma )_{[h]},$$
the subscript $`[h]`$ indicating to dilate the coefficients of the operator in the form, $`f_h(s,u):=f(s,h^{\frac{1}{2}}u).`$
Expanding the coefficients in Taylor series at $`h=0`$, we obtain the asymptotic expansion
$`(2.12)`$
$$\mathrm{\Delta }_h\underset{m=0}{\overset{\mathrm{}}{}}h^{(2+m/2)}_{2m/2}$$
where $`_2=L^2,`$ $`_{3/2}=0`$ and where
$`(2.13).`$
$$_1=2L^1[i\frac{}{s}+\frac{1}{2}\{\underset{j=1}{\overset{n}{}}_{u_j}^2\underset{ij=1}{\overset{n}{}}K_{ij}(s)u_iu_j\}]$$
We observe that $`_1`$ is $`2L^1`$ times the distinguished element $``$ of the adapted model, and that $``$ has weight $`1`$ when the $`u`$-variables are given their natural weights $`1/2`$.
As discussed in §1. 4-5, it will be helpful to rescale the variables once again to make them weightless. Hence we change variables to $`x=L^{\frac{1}{2}}u`$ and rewrite $`\mathrm{\Delta }_h`$ and the $`_{2\frac{n}{2}}`$’s in terms of the $`x`$-variables. For instance we will henceforth write $``$ in the form:
$$=i\frac{}{s}+\frac{1}{2}[\underset{j=1}{\overset{n}{}}L^1_{x_j}^2\underset{ij=1}{\overset{n}{}}LK_{ij}(s)x_ix_j].$$
We further note that the operators $`_{2m/2}`$ now satisfy the periodicity condition (1.3.8ii) : Indeed, as noted above, $`\mathrm{\Delta }`$ has this property when expressed in normal coordinates, and the various transversal rescalings and the conjugations by $`T_h,M_h`$ preserve it. To indicate that an operator has this periodicity property and hence acts on $`_T`$, we will subscript the appropriate spaces of operators with a “T.”
From (2.11) we see that the terms in $`_{2\frac{m}{2}}`$ are of the form
$$\begin{array}{ccc}x^m\hfill & x^{m2}\hfill & x^{m4}\hfill \\ x^{m2}D_x^2\hfill & x^{m3}D_x\hfill & \\ x^{m2}D_s\hfill & x^{m4}D_s\hfill & x^{m4}D_s^2\hfill \end{array}$$
hence
$`(2.14a)`$
$$_{2\frac{m}{2}}\mathrm{\Psi }^2(S^1)_T_ϵ^{m4}+\mathrm{\Psi }^1(S^1)_T_ϵ^{m2}+\mathrm{\Psi }^o(S^1)_T\times _ϵ^m$$
where the subscript $`ϵ`$ indicates that the Weyl symbol is a polynomial with the parity of $`m`$. Moreover, in the $`s`$ variable, $`_{2\frac{m}{2}}`$ is a polynomial in $`D_s`$ of degree at most two, so we can refine (2.14a) to the statement
$`(2.14b)`$
$$_{2\frac{m}{2}}C_T^{\mathrm{}}(S_L^1,_ϵ^{m4})D_s^2+C_T^{\mathrm{}}(S_L^1,_ϵ^{m2})D_s+C_T^{\mathrm{}}(S_L^1,_ϵ^m)$$
Comparing with (1.3.4) we see that
$$=\mu (𝒜_L^{})D_s\mu (𝒜_L^{})^1$$
where $`𝒜_L`$ is the weightless Wronskian matrix $`D_{L^{\frac{1}{2}}}a^{}`$, that is
$$𝒜_L:=\left(\begin{array}{cc}L^{\frac{1}{2}}Im\dot{Y}\hfill & L^{\frac{1}{2}}Re\dot{Y}\hfill \\ L^{\frac{1}{2}}ImY\hfill & L^{\frac{1}{2}}ReY\hfill \end{array}\right).$$
This motivates the conjugation of (2.11) to the (untwisted) model. We therefore put
$$𝒟_h=\mu (𝒜_L^{})^1\mathrm{\Delta }_h\mu (𝒜_L^{})$$
which has the asymptotic expansion
$`(2.15)`$
$$𝒟_h\underset{m=o}{\overset{\mathrm{}}{}}h^{(2+\frac{m}{2})}𝒟_{2\frac{m}{2}}$$
with $`𝒟_2=I,𝒟_{\frac{3}{2}}=0,𝒟_1=D_s`$. Conjugation by $`\mu (𝒜_L^{})`$ preserves weights, homogeneity and parity in the variables $`(x,D_x)`$. It also transforms $`D_s`$ into $`D_s`$ plus a term quadratic in $`(x,D_x)`$. Hence we find easily that $`𝒟_{2\frac{m}{2}}`$ has weight -1 for each m and that
$`(2.16a)`$
$$𝒟_{2\frac{m}{2}}\mathrm{\Psi }^2()_ϵ^{m4}+\mathrm{\Psi }^1(^)_ϵ^{m2}+\mathrm{\Psi }^o()_ϵ^m$$
or, analogously to (2.14b),
$`(2.16b)`$
$$𝒟_{2\frac{m}{2}}C^{\mathrm{}}(,_ϵ^{m4})D_s^2+C^{\mathrm{}}(,_ϵ^{m2})D_s+C^{\mathrm{}}(,_ϵ^m)$$
Of course, conjugation with $`\mu (𝒜_L^{})`$ also alters the periodicity property of the terms $`𝒟_{2\frac{m}{2}}`$. From (1.3.7a) and (1.3.8) we see in fact that they transform like operators in the twisted model, i.e. on the quantum mapping cylinder of $`r_\alpha (L).`$ More precisely, we have
$`(2.17a)`$
$$𝒟_{2\frac{m}{2}}|_{s+L}=\mu (𝒜_L^{})^1(_{2\frac{m}{2}}\mu (𝒜_L^{})|_{(s+L)})=\mu (r_\alpha (L))\mu (𝒜_L^{})^1\mu (T)^{}_{2\frac{m}{2}}|_{(s+L)}\mu (T)\mu (𝒜_L^{})\mu (r_\alpha (L))^{}$$
$$=\mu (r_\alpha (L))\mu (𝒜_L^{})^{}_{2\frac{m}{2}}\mu (𝒜_L^{})\mu (r_\alpha (L)^{}=\mu (r_\alpha (L))𝒟_{2\frac{m}{2}}\mu (r_\alpha (L))^{}.$$
Equivalently, in terms of the matrix elements in the basis of Hermite functions, we have
$`(2.17b).`$
$$𝒟_j\gamma _q,\gamma _r(s+L)=e^{i(\kappa _q\kappa _r)s}𝒟_j\gamma _q,\gamma _r(s)$$
To indicate that these operators act on $`_\alpha `$ we henceforth subscript the appropriate spaces of operators with an $`\alpha `$.
To render these terms periodic in $`s`$, we have to conjugate further to the (untwisted) model under $`\mu (r_\alpha )`$. We record the resulting expressions, since they will be used later on. In the notation $`\stackrel{~}{𝒜}_L^{}(s):=𝒜_L^{}(s)r_\alpha (s)`$ the principal operator becomes, by Proposition (1.3.8),
$`(2.18)`$
$$\mu (\stackrel{~}{𝒜}_L^{})\mu (\stackrel{~}{𝒜})^1=$$
Since conjugation by $`\mu (r_\alpha )`$ also preserves weight, homogeneity and parity, we further have:
$`(2.19)`$
$$_h:=\mu (\stackrel{~}{a})\mathrm{\Delta }_h\mu (\stackrel{~}{a})^{}\underset{m=o}{\overset{\mathrm{}}{}}h^{(2+\frac{m}{2})}_{2\frac{m}{2}}$$
with $`_2=I,_{\frac{3}{2}}=0,_1=`$ and with all coefficients of weight -2 and periodic, that is,
$`(2.20a)`$
$$_{2\frac{m}{2}}\mathrm{\Psi }^2(S^1)_ϵ^{m4}+\mathrm{\Psi }^1(S^1)_ϵ^{m2}+\mathrm{\Psi }^o(S^1)_ϵ^m$$
or, analogously to (2.14b),
$`(2.20b).`$
$$_{2\frac{m}{2}}C^{\mathrm{}}(S_L^1,_ϵ^{m4})^2+C^{\mathrm{}}(S_L^1,_ϵ^{m2})+C^{\mathrm{}}(S_L^1,_ϵ^m)$$
Our aim is now to put $`\mathrm{\Delta }`$, or $`\mathrm{\Delta }_h`$ for certain $`h`$, into a semi-classical normal form. This normal form will be, at first, only a formal normal form for $`\mathrm{\Delta }_h`$ as a formal $`h`$pseudo-differential operator. Here, a formal $`h`$-pseudodifferential operator of order m on $`^N`$ is the Weyl quantization
$`(2.21a)`$
$$a^w(x,hD;h)u(x)=(2\pi h)^ne^{\frac{i}{h}xy,\xi }a(\frac{1}{2}(x+y),\xi ;h)u(y)𝑑y𝑑\xi $$
of an amplitude $`a`$ belonging to the space $`S_{cl}^m(T^{}^N)`$ of asymptotic sums
$`(2.21b)`$
$$a(x,\xi ,h)h^m\underset{j+o}{\overset{\mathrm{}}{}}a_j(x,\xi )h^j$$
with $`a_jC^{\mathrm{}}(R^{2N}).`$ Such operators form an algebra $`\mathrm{\Psi }_h^{}(^N)`$ under composition, with $`\mathrm{\Psi }_h^m`$ the subspace of mth order elements. (See \[Sj\] for further background and references). We will also be concerned with the slightly different situation of $`h`$\- pseudodifferential operators on $`S^1\times ^n`$, which are defined similarly using local coordinates. Combining the $`h`$ filtration with the previous filtrations of $`\mathrm{\Psi }^{}()𝒲^{}(^n)`$ we get the triply filtered algebra
$$\mathrm{\Psi }_h^{(,,)}(\times ^n)$$
$$\mathrm{\Psi }_h^{(k,m,\frac{n}{2})}:=h^k\mathrm{\Psi }^m()𝒲^{\frac{n}{2}},$$
with $`k+m+\frac{n}{2}`$ the total order of an element (and similarly with $`S^1`$ replacing $``$.)
The following lemma will prepare for the normal form. We state it in terms of the $``$-operators of the model since the periodicity properties are simplest there. The notation “$`A|_o`$” will be used for the restriction of an operator $`A\mathrm{\Psi }^{}(S^1\times ^n)`$ of the model to elements of $``$-weight zero. Equivalently, after conjugation by $`\mu (r_\alpha )`$ to the twisted model, to elements of $`D_s`$-weight zero , that is, to functions independent of $`s`$ in $`_\alpha `$. If we write the latter $`A`$ in the form $`A_2D_s^2+A_1D_s+A_o`$, then $`A|_o`$ is $`A_o|_o`$ (restricted to weight 0 elements). We have:
(2.22)Lemma There exists an $`L`$-dependent $`h`$-pseudodifferential operator $`W_h=W_h(s,x,D_x)`$ on $`L^2(S_L^1\times ^n)`$ such that, for each $`sS_L^1`$,
$$W_h(s,x,D_x):L^2(^n)L^2(^n)$$
is unitary, and such that
$$W_h^{}_hW_hh^2L^2+2h^1L^1+\underset{j=0}{\overset{\mathrm{}}{}}h^{\frac{j}{2}}_{2\frac{j}{2}}^{\mathrm{}}(s,D_s,y,D_y)$$
where
(i) $`_{2\frac{j}{2}}^{\mathrm{}}(s,D_s,x,D_x)=_{2\frac{j}{2}}^{\mathrm{},2}^2+_{2\frac{j}{2}}^{\mathrm{},1}+_{2\frac{j}{2}}^{\mathrm{},o},`$ with $`_{2\frac{j}{2}}^{\mathrm{},k}C^{\mathrm{}}(S_L^1,_ϵ^{j2k});`$
(ii) $`_{2j}^{\mathrm{}}(s,D_s,x,D_x)|_o=_{2j}^{\mathrm{},o}(s,x,D_x)|_o=f_j(I_1,\mathrm{},I_n)|_o`$ for certain polynomials $`f_j`$ of degree j+2 on $`^n,`$ i.e. $`f_j(I_1,\mathrm{},I_n)𝒫_{}^{j+2}`$
(iii) $`_{2\frac{2k+1}{2}}^{\mathrm{}}(s,D_s,x,D_x)|_o=_{2\frac{2k+1}{2}}^{\mathrm{},o}(s,x,D_x)|_o=0;`$
(iv) The harmonic oscillators are weightless and all of the $``$’s have weight -2.
Proof: The operator $`W_h`$ will be constructed as the asymptotic product
$`(2.23)`$
$$W_h:=\mu (r_\alpha )^{}\mathrm{\Pi }_{k=1}^{\mathrm{}}W_{h\frac{k}{2}}\mu (r_\alpha )$$
of weightless unitary $`h`$-pseudodifferential operators on $`^n`$, with
$`(2.23a)`$
$$W_{h\frac{k}{2}}:=exp(ih^{\frac{k}{2}}Q_{\frac{k}{2}})$$
and with $`h^{\frac{k}{2}}Q_{\frac{k}{2}}h^{\frac{k}{2}}^{\mathrm{}}(S_L^1)^{k+2}`$ of total order 1. The product will converge, for each s, to a unitary operator in $`\mathrm{\Psi }_h^o(^n)`$ (see \[Sj\] for discussion of asymptotic products).
To see what is involved, we first construct a weightless $`Q_{\frac{1}{2}}(s,x,D_x)C^{\mathrm{}}(S_L^1)_ϵ^3`$ such that
$`(2.24a)`$
$$e^{ih^{\frac{1}{2}}Q_{\frac{1}{2}}}_he^{ih^{\frac{1}{2}}Q_{\frac{1}{2}}}|_o=[h^2L^2+2h^1L^1+_o^{\frac{1}{2}}+\mathrm{}]|_o$$
where the dots $`\mathrm{}`$ indicate higher powers in $`h`$. The operator $`Q_{\frac{1}{2}}`$ then must satisfy the commutation relation
$`(2.24b)`$
$$\{[L^1,Q_{\frac{1}{2}}]+_{\frac{1}{2}}\}|_o=0.$$
To solve for $`Q_{\frac{1}{2}}`$, it is convenient to conjugate back to the $`𝒟_{2\frac{m}{2}}`$’s of the twisted model by $`\mu (r_\alpha )`$, since this transforms $``$ into $`D_s`$. The commutation relation thus becomes
$`(2.24c)`$
$$\{[L^1D_s,\mu (r_\alpha )^{}Q_{\frac{1}{2}}\mu (r_\alpha )]+𝒟_{\frac{1}{2}}\}|_o=0,$$
that is,
$`(2.24d)`$
$$L^1_s\{\mu (r_\alpha )^{}Q_{\frac{1}{2}}\mu (r_\alpha )\}|_o=i\{𝒟_{\frac{1}{2}}\}|_o$$
where $`_sA`$ is the Weyl operator whose complete symbol is the $`s`$-derivative of that of $`A`$. Since (2.24d) is simpler than (2.24b), we henceforth conjugate everything by $`\mu (r_\alpha )`$, and relabel the operators $`\mu (r_\alpha )^{}Q\mu (r_\alpha )`$ by $`\stackrel{~}{Q}.`$ The resulting $`𝒟`$’s then have the twisted model periodicity properties (2.17a-b). Our problem is then to solve (2.24d) with an operator $`\stackrel{~}{Q}_{\frac{1}{2}}`$ satisfying (2.17a-b).
To do so, we first observe that under conjugation by $`\mu (r_\alpha )`$, elements of $``$-weight zero transform to elements of $`D_s`$\- weight zero, and hence it suffices to solve for the matrix elements $`\stackrel{~}{Q}_{\frac{1}{2}}\gamma _q,\gamma _r`$. It follows from (2.17b) that the solution is unique for $`rq`$, while for $`r=q`$ the matrix element $`\stackrel{~}{Q}_{\frac{1}{2}}\gamma _q,\gamma _q`$ is periodic and hence a necessary and sufficient condition for solvability by an operator with the correct periodicity is that, for all $`q𝐍^n,`$
$`(2.25).`$
$$_o^L𝒟_{\frac{1}{2}}\gamma _q,\gamma _q𝑑s=0$$
In view of (2.16a-b), we have
$$𝒟_{\frac{1}{2}}\mathrm{\Psi }^1()_\alpha _ϵ^1+\mathrm{\Psi }^o()_\alpha _ϵ^3,$$
in fact a simple calculation shows that $`𝒟_{\frac{1}{2}}\mathrm{\Psi }^o()_\alpha _ϵ^3`$. (Recall that the subscript $`\alpha `$ indicates the periodicity property (2.17a-b).)
We now observe that if $`A\mathrm{\Psi }^{}()_ϵ^m`$ then
$`(2.26)`$
$$A:\left\{\begin{array}{cc}L_\pm ^2L_\pm ^2\hfill & m\text{even}\hfill \\ :L_\pm ^2L_{}^2\hfill & m\text{odd}\hfill \end{array}\right\},$$
where $`L_+^2`$ (resp. $`L_{}^2`$) denotes the subspace of $`L^2()L^2(^n)`$ spanned by even (resp. odd) functions of $`y^n`$ with coefficients in $`s`$.
It follows that
$$𝒟_{\frac{1}{2}}:L_+^2L_{}^2,$$
This parity reversing property implies the vanishing of the diagonal matrix elements in the basis $`\{\gamma _q\},`$ hence (2.25) does hold. A unique solution of (2.24d) is specified by the condition that
$$_o^{2\pi }\stackrel{~}{Q}_{\frac{1}{2}}\gamma _q,\gamma _q𝑑s=0.$$
To solve the equation (2.24d) we rewrite it in terms of complete Weyl symbols. We will use the notation $`A(s,x,\xi )`$ for the complete Weyl symbol of the operator $`A(s,x,D_x)`$. Then (2.24d) becomes
$`(2.27a)`$
$$L^1_s\stackrel{~}{Q}_{\frac{1}{2}}(s,x,\xi )=i𝒟_{\frac{1}{2}}|_o(s,x,\xi )$$
with
$$\stackrel{~}{Q}_{\frac{1}{2}}(s+L,x,\xi )=\stackrel{~}{Q}_{\frac{1}{2}}(s,r_\alpha (L)(x,\xi )).$$
We solve (2.27a) with the Weyl symbol
$$\stackrel{~}{Q}_{\frac{1}{2}}(s,x,\xi )=\stackrel{~}{Q}_{\frac{1}{2}}(0,x,\xi )+L_0^si𝒟_{\frac{1}{2}}|_o(u,x,\xi )du$$
where $`\stackrel{~}{Q}_{\frac{1}{2}}(0,x,\xi )`$ is determined by the consistency condition
$`(2.27b)`$
$$\stackrel{~}{Q}_{\frac{1}{2}}(L,x,\xi )\stackrel{~}{Q}_{\frac{1}{2}}(0,x,\xi )=L_0^Li𝒟_{\frac{1}{2}}|_o(u,x,\xi )du$$
or in view of the periodicity condition in (2.27a),
$`(2.27c)`$
$$\stackrel{~}{Q}_{\frac{1}{2}}(0,r_\alpha (x,\xi ))\stackrel{~}{Q}_{\frac{1}{2}}(0,x,\xi )=L_0^Li𝒟_{\frac{1}{2}}|_o(u,x,\xi )du.$$
To solve, we use that $`𝒟_{\frac{1}{2}}|_o(u,x,\xi )`$ is a polynomial of degree 3 in $`(x,\xi )`$. Also, as is customary in such calculations, we switch to complex coordinates $`z_j=x_j+i\xi _j`$ and $`\overline{z}_j=x_ji\xi _j`$ in which the action of $`r_\alpha (L)`$ is diagonal. With a little abuse of notation, we will continue to denote the Weyl symbols, qua functions of the $`z_j,\overline{z}_j`$’s by their previous expressions. We also suppress the subscripts by using vector notation $`z,\overline{z}`$ and $`e^{i\alpha }`$. Thus, (2.27c) becomes
$`(2.28)`$
$$\stackrel{~}{Q}_{\frac{1}{2}}(0,e^{i\alpha }z,e^{i\alpha }\overline{z})\stackrel{~}{Q}_{\frac{1}{2}}(0,z,\overline{z})=L_0^Li𝒟_{\frac{1}{2}}|_o(u,z,\overline{z})du.$$
We now use that $`𝒟_{\frac{1}{2}}(u,z,\overline{z})`$ is a polynomial of degree 3 to solve (2.27c). If we put
$$\stackrel{~}{Q}_{\frac{1}{2}}(s,z,\overline{z})=\underset{|m|+|n|3}{}q_{\frac{1}{2};mn}(s)z^m\overline{z}^n$$
and
$$𝒟_{\frac{1}{2}}|_o(s,z,\overline{z})du=\underset{|m|+|n|3}{}d_{\frac{1}{2};mn}(s)z^m\overline{z}^n$$
then (2.28) becomes
$`(2.29).`$
$$\underset{|m|+|n|3}{}(1e^{(mn)\alpha })q_{\frac{1}{2};mn}(0)z^m\overline{z}^n=iL^2\underset{|m|+|n|3}{}\overline{d}_{\frac{1}{2};mn}z^m\overline{z}^n$$
Since there are no terms with $`m=n`$ in this (odd-index) equation, and since the $`\alpha _j`$’s are independent of $`\pi `$ over $``$, there is no obstruction to the solution of (2.29).
For simplicity of notation we will express the solution in the form
$`(2.30)`$
$$\stackrel{~}{Q}_{\frac{1}{2}}|_o=iL𝒟_{\frac{1}{2}}(s)𝑑s|_o$$
where $``$ denotes the integration procedure just defined, that is, the indefinite integral satisfying (2.27b). Since the integration only involves the $`s`$-coefficients, the solution is a pseudodifferential operator on $`^n`$ with the same order, same order of vanishing, and same parity as the restriction of $`𝒟_{\frac{1}{2}}`$ to elements of weight zero. We then extend it all of to $`_\alpha `$ as a pseudodifferential operator with the same properties by stipulating that
$`(2.31a)`$
$$\stackrel{~}{Q}_{\frac{1}{2}}M_h=M_h\stackrel{~}{Q}_{\frac{1}{2}}.$$
Then, as desired,
$`(2.31b)`$
$$\stackrel{~}{Q}_{\frac{1}{2}}\mathrm{\Psi }^o(^1)_ϵ^3.$$
The conjugate by $`\mu (r_\alpha )`$ then defines a periodic operator satisfying (2.24b) and hence a unitary element $`W_{h\frac{1}{2}}\mathrm{\Psi }_h^o(S^1\times ^n)`$ satisfying (2.24a). The twisted unitary operator with exponent $`\stackrel{~}{Q}_{\frac{1}{2}}`$, i.e.the image of $`W_{h\frac{1}{2}}`$ under conjugation by $`\mu (r_\alpha )`$, will be decorated with a tilde, $`\stackrel{~}{W}_{h\frac{1}{2}}.`$
Since $`h^{\frac{1}{2}}\stackrel{~}{Q}_{\frac{1}{2}}`$ is of total order 1, $`h^{\frac{1}{2}}ad(\stackrel{~}{Q}_{\frac{1}{2}})`$ (with $`ad(A)B:=[B,A]`$) preserves the total order in $`\mathrm{\Psi }_h^{(,,)}`$, and hence $`\stackrel{~}{W}_{h\frac{1}{2}}`$ is an order-preserving automorphism. It is moreover independent of $`D_s`$ and has an odd polynomial Weyl symbol, so that
$`(2.32)`$
$$h^{\frac{1}{2}}ad(\stackrel{~}{Q}_{\frac{1}{2}}):h^{\frac{k}{2}}\mathrm{\Psi }^l()_ϵ^mh^{\frac{k+1}{2}}[\mathrm{\Psi }^{l1}()_ϵ^{m+3}+\mathrm{\Psi }^l()_ϵ^{m+1}].$$
Finally, since $`𝒟_{\frac{1}{2}}`$ has weight -2, $`L`$ has degree 1 and since $`_o^s()𝑑s`$ has degree 1 in the metric scaling, we see that $`\stackrel{~}{Q}_{\frac{1}{2}}`$ is weightless. Alternatively, the $`d_{\frac{1}{2};m,n}`$’s have weight -2,the variables $`z`$ have weight 0 and hence the $`q_{\frac{1}{2};m,n}`$’s have weight 0.
Consider now the element
$$𝒟_h^{\frac{1}{2}}:=\stackrel{~}{W}_{h\frac{1}{2}}^{}𝒟_h\stackrel{~}{W}_{h\frac{1}{2}}\mathrm{\Psi }_h^2(^1\times ^n).$$
Using only the $`h`$-filtration, we expand
$`(2.33)`$
$$𝒟_h^{\frac{1}{2}}\underset{n=o}{\overset{\mathrm{}}{}}h^{2+\frac{n}{2}}\underset{j+m=n}{}\frac{i^j}{j!}(ad\stackrel{~}{Q}_{\frac{1}{2}})^j𝒟_{2\frac{m}{2}}$$
$$:=h^2L^2+h^1L^1D_s+\underset{n=3}{\overset{\mathrm{}}{}}h^{2+\frac{n}{2}}𝒟_{2\frac{n}{2}}^{\frac{1}{2}}.$$
An obvious induction using (2.32) and (2.16a-b) gives that
$$ad(\stackrel{~}{Q}_{\frac{1}{2}})^j𝒟_{2\frac{m}{2}}C^{\mathrm{}}(,_ϵ^{m+j4})D_s^2+C^{\mathrm{}}(,_ϵ^{m+j2})D_s+C^{\mathrm{}}(,_ϵ^{m+j}).$$
It follows that $`𝒟_{2\frac{n}{2}}^{\frac{1}{2}}`$ has the same filtered structure (2.16b) as $`𝒟_{2\frac{n}{2}}.`$
We carry this procedure out one more step before arguing inductively, since the even steps behave differently from the odd ones. We thus seek an element $`\stackrel{~}{Q}_1(s,x,D_x)\mathrm{\Psi }^{}(S^1\times ^n)`$ and an element $`f_o(I_1,\mathrm{},I_n)𝒜`$ so that
$$𝒟_h^1:=\stackrel{~}{W}_{h1}^{}𝒟^{\frac{1}{2}}\stackrel{~}{W}_{h1}=h^2L^2+h^1L^1D_s+h^{\frac{1}{2}}𝒟_{\frac{1}{2}}^{\frac{1}{2}}+𝒟_o^1(s,D_s,x,D_x)+\mathrm{}$$
with
$`(2.34a)`$
$$𝒟_o^1(s,D_s,x,D_x)|_o=f_o(I_1,\mathrm{},I_n)$$
with $`\stackrel{~}{W}_{h1}=e^{ih\stackrel{~}{Q}_1},`$ and where the dots signify terms of higher order in $`h`$. Note that $`𝒟_{\frac{1}{2}}^1=𝒟_{\frac{1}{2}}^{\frac{1}{2}}`$, so that (2.33a) would imply
$`(2.34b)`$
$$\{h^{\frac{1}{2}}𝒟_{\frac{1}{2}}^1+𝒟_o^1\}|_o=f_o(I_1,\mathrm{},I_n).$$
The condition on $`\stackrel{~}{Q}_1`$ is then
$`(2.35a)`$
$$\{[D_s,\stackrel{~}{Q}_1]+𝒟_o^{\frac{1}{2}}\}|_o=f_o(I_1,\mathrm{},I_n)$$
or equivalently
$`(2.35b).`$
$$_s\stackrel{~}{Q}_1|_o=\{𝒟_o^{\frac{1}{2}}+f_o(I_1,\mathrm{},I_n)\}|_o$$
As above, we first consider the solvability of (2.35b) one matrix element at a time. The off-diagonal matrix element equations again have a unique solution, but because $`𝒟_o^{\frac{1}{2}}`$ is even there is now a condition on the solvability, with a periodic solution, of the diagonal ones:
$`(2.36)`$
$$\frac{1}{L}_o^L𝒟_o^{\frac{1}{2}}\gamma _q,\gamma _q𝑑s=f_o(q_1+\frac{1}{2},\mathrm{},q_n+\frac{1}{2}).$$
These conditions (2.36) determine $`f_o`$ and hence the operator $`f_o(I_1,\mathrm{},I_n).`$ To analyse the properties of $`f_o`$ and of $`\stackrel{~}{Q}_1`$, we use (1.1b), (2.16b) and (2.31c) to conclude that
$$𝒟_o^{\frac{1}{2}}=𝒟_o+ad\stackrel{~}{Q}_{\frac{1}{2}}(ad\stackrel{~}{Q}_{\frac{1}{2}}D_s+𝒟_{\frac{1}{2}})$$
$$C^{\mathrm{}}(,_ϵ^o)D_s^2+C^{\mathrm{}}(,_ϵ^2)D_s+C^{\mathrm{}}(,_ϵ^4).$$
We also observe that (2.32) is equivalent to
$`(2.37)`$
$$f_o(I_1,\mathrm{},I_n)=\frac{1}{L}_o^L_{T^n}V_t^{}𝒟_o^{\frac{1}{2}}|_oV_tdtds$$
where $`T^n=^n/^n`$ is the n-torus, where $`tT^n`$ and where
$$V(t_1,\mathrm{},t_n):=expit_1I_1\mathrm{}expit_nI_n.$$
Since $`V_t`$ belongs to the metaplectic representation, we have by metaplectic covariance of the Weyl calculus that
$`(2.38).`$
$$f_o(I_1,\mathrm{},I_n)𝒫_{}^2^4$$
We then can solve for $`\stackrel{~}{Q}_1`$ in the form
$`(2.39)`$
$$\stackrel{~}{Q}_1(s,y,D_y):=iL[\{𝒟_o^{\frac{1}{2}}f_o(I_1,\mathrm{},I_n)\}|_ods]\mathrm{\Psi }^o(^1)_ϵ^4$$
where $``$ is the indefinite integration in $`s`$ with zero mean value and where the operator on the right is extended to $`L^2(^1\times ^n)`$ using (2.35b).
To make the integration more concrete and to analyze the solution, we again express everything in terms of complete Weyl symbols. Thus we rewrite (2.35b) in the form
$`(2.40a)`$
$$L^1_s\stackrel{~}{Q}_1(s,z,\overline{z})=i\{𝒟_o^{\frac{1}{2}}|_o(s,z,\overline{z})f_o(|z_1|^2,\mathrm{},|z_n|^2)\}$$
or equivalently
$`(2.40b)`$
$$\stackrel{~}{Q}_1(s,z,\overline{z})=\stackrel{~}{Q}_1(0,z,\overline{z})iL_0^s[𝒟_o^{\frac{1}{2}}|_o(u,z,\overline{z})f_o(|z_1|^2,\mathrm{},|z_n|^2)]𝑑u$$
and solve simeltaneously for $`\stackrel{~}{Q}_1`$ and $`f_o`$. The consistency condition determining a unique solution is that
$`(2.41a)`$
$$\stackrel{~}{Q}_1(L,z,\overline{z})=\stackrel{~}{Q}_1(0,z,\overline{z})iL_0^L[𝒟_o^{\frac{1}{2}}|_o(u,z,\overline{z})f_o(|z_1|^2,\mathrm{},|z_n|^2)]𝑑u.$$
or in view of the twisted periodicity condition
$`(2.41b)`$
$$\stackrel{~}{Q}_1(0,e^{i\alpha }z,e^{i\alpha }\overline{z})\stackrel{~}{Q}_1(0,z,\overline{z})=iL\{_0^L𝒟_o^{\frac{1}{2}}|_o(u,z,\overline{z})duLf_o(|z_1|^2,\mathrm{},|z_n|^2)\}.$$
As before we use that $`𝒟_o^{\frac{1}{2}}|_o(u,z,\overline{z})`$ is a polynomial of degree 4 to solve the equation. We put
$`(2.42)`$
$$\stackrel{~}{Q}_1(s,z,\overline{z})=\underset{|m|+|n|4}{}q_{1;mn}(s)z^m\overline{z}^n,f_o(|z_1|^2,\mathrm{},|z_n|^2)=\underset{|k|2}{}c_{ok}|z|^{2k}$$
and
$$𝒟_o^{\frac{1}{2}}|_o(s,z,\overline{z})du:=\underset{|m|+|n|4}{}d_{o;mn}^{\frac{1}{2}}(s)z^m\overline{z}^n,\overline{d}_{o;mn}^{\frac{1}{2}}:=\frac{1}{L}_o^Ld_{o;mn}^{\frac{1}{2}}(s)𝑑s$$
As above, we can solve for the off-diagonal coefficients,
$`(2.43a)`$
$$q_{1;mn}(0)=iL^2(1e^{i(mn)\alpha })^1\overline{d}_{o;mn}^{\frac{1}{2}}$$
and must set the diagonal coefficients equal to zero. The coefficients $`c_{ok}`$ are then determined by
$`(2.43b)`$
$$c_{ok}=\overline{d}_{1;kk}^{\frac{1}{2}}.$$
It is evident that $`\stackrel{~}{Q}_1`$ and $`f_o(I_1,\mathrm{},I_n)`$ are even polynomial pseudodifferential operators of degree 4 in the variables $`(x,D_x)`$, that $`\stackrel{~}{Q}_1`$ is weightless under metric rescalings and that the coefficients $`c_{ok}`$ are of weight -2. They are essentially the same as the residual QBNF invariants.
We now proceed inductively to define self-adjoint polynomial pseudodifferential operators $`\stackrel{~}{Q}_{\frac{j}{2}},f_j(I_1,\mathrm{},I_n)`$, unitary $`h`$-pseudodifferential operators $`\stackrel{~}{W}_{h\frac{j}{2}}:=exp(ih^{\frac{j}{2}}\stackrel{~}{Q}_{\frac{j}{2}})`$ and approximate semi-classical normal forms $`𝒟_h^{\frac{k}{2}}`$ such that:
(2.44)
$$\begin{array}{c}(i)𝒟_h^{\frac{k}{2}}:=\stackrel{~}{W}_{h\frac{k}{2}}^{}\stackrel{~}{W}_{h\frac{k1}{2}}^{}\mathrm{}\stackrel{~}{W}_{h\frac{1}{2}}^{}𝒟_h\stackrel{~}{W}_{h\frac{1}{2}}\mathrm{}\stackrel{~}{W}_{h\frac{k1}{2}}\stackrel{~}{W}_{h\frac{k}{2}};\hfill \\ (ii)𝒟_h^{\frac{k}{2}}h^2+h^1D_s+_{n=3}^{\mathrm{}}h^{2+\frac{n}{2}}𝒟_{2\frac{n}{2}}^{\frac{k}{2}};\hfill \\ (iii)\text{ for}kn2,𝒟_{2\frac{n}{2}}^{\frac{k}{2}}=𝒟_{2\frac{n}{2}}^{\frac{n2}{2}}=[D_s,\stackrel{~}{Q}_{\frac{n}{2}1}]+𝒟_{\frac{n}{2}}^{\frac{n3}{2}}\hfill \\ (iv)\text{for}kn2,𝒟_{2\frac{n}{2}}^{\frac{k}{2}}|_o\left\{\begin{array}{cc}=f_j(I_1,\mathrm{},I_n)\hfill & n=2j\hfill \\ =0\hfill & n=2j+1\hfill \end{array}\right\}\hfill \\ (v)𝒟_{2\frac{m}{2}}^{\frac{k}{2}}C^{\mathrm{}}(,_ϵ^{m4})D_s^2+C^{\mathrm{}}(,_ϵ^{m2})D_s+C^{\mathrm{}}(,_ϵ^m);\hfill \\ (vi)\stackrel{~}{Q}_{\frac{k}{2}}C^{\mathrm{}}(,_ϵ^{k+2});\hfill \\ (vii)\text{the conjugates under}\mu (r_\alpha )\text{of the above operators are periodic of period L in s};\hfill \\ (viii)f_j(I_1,\mathrm{},I_n)𝒫_{}^{j+2};(ix)wgt(f_j)=2.\hfill \end{array}$$
To check the induction, let us assume these properties hold for $`kN`$. Then
$`(2.45)`$
$$𝒟_h^{\frac{N+1}{2}}=\stackrel{~}{W}_{h\frac{N+1}{2}}^{}𝒟_h^{\frac{N}{2}}\stackrel{~}{W}_{h\frac{N+1}{2}}=$$
$$\begin{array}{c}=h^2L^2+h^1L^1D_s+_{n=3}^{N+2}h^{2+\frac{n}{2}}𝒟_{2\frac{n}{2}}^{\frac{N}{2}}+h^{\frac{N1}{2}}\{[L^1D_s,\stackrel{~}{Q}_{\frac{N+1}{2}}]+𝒟_{\frac{N1}{2}}^{\frac{N}{2}}\}+_{n=N+4}^{\mathrm{}}h^{2+\frac{n}{2}}𝒟_{2\frac{n}{2}}^{\frac{N}{2}}\hfill \\ +_{m=2}^{\mathrm{}}h^{m\frac{N+1}{2}}\frac{i^m}{m!}(ad\stackrel{~}{Q}_{\frac{N+1}{2}})^m(D_s)+_{m=1}^{\mathrm{}}_{p=3}^{\mathrm{}}h^{m\frac{N+1}{2}+\frac{p}{2}2}\frac{i^m}{m!}(ad\stackrel{~}{Q}_{\frac{N+1}{2}})^m(𝒟_{2\frac{p}{2}}^{\frac{N}{2}}).\hfill \end{array}$$
We note that
$$𝒟_{2\frac{N+3}{2}}^{\frac{N+1}{2}}=[L^1D_s,\stackrel{~}{Q}_{\frac{N+1}{2}}]+𝒟_{\frac{N1}{2}}^{\frac{N}{2}}$$
and that further conjugations will not alter this (or lower order) term(s), proving (iii).
When $`N+1`$ is odd, $`\stackrel{~}{Q}_{\frac{N+1}{2}}`$ must solve:
$`(2.46a)`$
$$\{[L^1D_s,\stackrel{~}{Q}_{\frac{N+1}{2}}]+𝒟_{\frac{N1}{2}}^{\frac{N}{2}}\}|_o=0$$
while if $`N+1`$ is even $`\stackrel{~}{Q}_{\frac{N+1}{2}}`$ and $`f_{\frac{N1}{2}}`$ must solve
$`(2.46b)`$
$$\{[L^1D_s,\stackrel{~}{Q}_{\frac{N+1}{2}}]+𝒟_{\frac{N1}{2}}^{\frac{N}{2}}\}|_o=f_{\frac{N1}{2}}(I_1,\mathrm{},I_n).$$
If $`N+1`$ is odd, (2.45v) implies that the diagonal matrix elements of $`𝒟_{\frac{N1}{2}}^{\frac{N}{2}}|_o`$ vanish. Hence, as in the case of $`\stackrel{~}{Q}_{\frac{1}{2}}`$, the solution of (2.46a) is given by
$`(2.47a)`$
$$\stackrel{~}{Q}_{\frac{N+1}{2}}|_o:=iL𝒟_{\frac{N1}{2}}^{\frac{N}{2}}|_ods$$
where the constant of integration $`\stackrel{~}{Q}_{\frac{N+1}{2}}|_o(0,z,\overline{z})`$ is defined so that the solution satisfies the periodicity condition analogous to (2.28). In the even case, as in the case of $`\stackrel{~}{Q}_1`$ we set
$`(2.47b)`$
$$f_{\frac{N1}{2}}(q_1+\frac{1}{2},\mathrm{},q_n+\frac{1}{2}):=\frac{1}{L}_o^L𝒟_{\frac{N1}{2}}^{\frac{N}{2}}\gamma _q,\gamma _q𝑑s,$$
or equivalently
$`(2.47c)`$
$$f_{\frac{N1}{2}}(I_1,\mathrm{},I_n):=\frac{1}{L}_o^L_{T^n}V_t^{}𝒟_{\frac{N1}{2}}^{\frac{N}{2}}|_oV_tdtds,$$
and
$`(2.47d)`$
$$\stackrel{~}{Q}_{\frac{N+1}{2}}|_o:=L\{𝒟_{\frac{N1}{2}}^{\frac{N}{2}}f_{\frac{N1}{2}}(I_1,\mathrm{},I_n)\}|_ods.$$
As above, $`\stackrel{~}{Q}_{\frac{N+1}{2}}|_o`$ is then extended to all of $`_\alpha `$ thru (2.31b). The indefinite integrations can be precisely defined, as above, by expressing everything as a polynomial in $`z,\overline{z}`$ and solving the resulting algebraic equations for the constant of integration. By (2.47 c-d) the solution has the parity of $`𝒟_{\frac{N1}{2}}^{\frac{N}{2}}`$, i.e. the parity of $`N1,`$ and by (2.45) the parity property propagates to the case $`N+1`$. The formula (2.47c) also shows that for even $`N1`$, $`f_{\frac{N1}{2}}(I_1,\mathrm{},I_n)`$ has the same order and weight properties as $`𝒟_{\frac{N1}{2}}^{\frac{N}{2}}`$. Since the latter may be written in the form $`aD_s^2+bD_s+c,`$ with $`cC^{\mathrm{}}(,_ϵ^{N+3})`$, the former lies in $`𝒫_{}^{}_ϵ^{N+3}`$, implying $`(viii)`$. The formula (2.47d) then implies $`(vi)`$. The periodicity property (vii) is maintained throughout.
It follows immediately that
$`(2.48)`$
$$\begin{array}{cc}𝒟_{2\frac{n}{2}}^{\mathrm{}}\hfill & =𝒟_{2\frac{n}{2}}^{\frac{n2}{2}}\hfill \\ & =[D_s,\stackrel{~}{Q}_{\frac{n}{2}1}]+𝒟_{2\frac{n}{2}}^{\frac{n3}{2}}.\hfill \end{array}$$
Conjugating back under $`\mu (r_\alpha )`$, we see that statement $`(i)`$ of the Lemma then follows from (2.48) and from (2.44iii-iv) and statement (ii) then follows from (2.44 viii). The weight property (ix) is visible from (2.47c). ∎
As a transitional step to the quantum Birkhoff normal form, let us rephrase the above Lemma in terms of the actual Fermi normal coordinates and inverse Planck constants $`\{r_{kq}\}`$ in place of $`(Lh)^1`$. We thus introduce the space $`𝒪^{}(N_\gamma ,\mathrm{\Gamma },\{r_{kq}\})`$ of semi-classical Hermite distributions associated to the non-homogenenous isotropic manifold $`\mathrm{\Gamma }:=^+\gamma `$, i.e.
$$\mathrm{\Gamma }:=\{(s,\sigma ,y,\eta )T^{}(N_\gamma ):\sigma =1,(y,\eta )=(0,0)\}.$$
It is, by definition, the union of the spaces $`𝒪^{m/2}`$ of elements of order $`\frac{m}{2},m𝐍`$, given in Fermi normal coordinates by asymptotic sums of the form
$$u(s,y,r_{kq}^1)(r_{kq})^{\frac{m}{2}}e^{ir_{kq}s}\underset{n=o}{\overset{\mathrm{}}{}}(r_{kq})^{\frac{n}{2}}f_n(s,\sqrt{r_{kq}}y)$$
with $`f_nC_q^{\mathrm{}}(,𝒮(^n))`$. Here, $`𝒮(^n)`$ is the Schwarz space, and the subscript q denotes the space of functions of this form satisfying
$$f(s+L,y)=e^{i\kappa _qs}f(s,ty).$$
Aside from the restriction to Schwartz functions and the half-integral orders, it is the image under the rescaling operator $`T_{(r_{kq})^1}`$ of the space $`𝒪^{}(N_\gamma ,\mathrm{\Lambda },\{r_{kq}\})`$ of semi-classical Lagrangean distributions associated to the non-homogenenous Lagrangean
$$\mathrm{\Lambda }:=graph(ds)=\{(s,\sigma ,y,\eta ):\sigma =1,\eta =0.\}.$$
Similar spaces of oscillatory functions could be, and implicitly are, defined in the model $`S^1\times ^n`$ but the scaling aspect comes out most naturally on $`N_\gamma `$. (For background on semi-classical Lagrangean distributions, see \[CV2\]).
Let us denote by $`𝒪_o^{}(N_\gamma ,\mathrm{\Lambda },\{r_{kq}\})`$ the subspace of elements of $`𝒪^{}(S^1\times ^n,\mathrm{\Lambda })`$ satisfying $`u(s,y,r_{kq}^1)0.`$ Let us also denote by $`W_{kq}`$ the transfer of $`W_{r_{kq}^1}`$ above to the normal bundle, i.e.
$$W_{kq}:=\mu (\stackrel{~}{a}_s)W_{r_{kq}^1}\mu (\stackrel{~}{a}_s)^{}.$$
For the sake of simplicity we will be a little negligent here of the rescalings. Under this operator, the kernel of $``$ goes over to the space $`\stackrel{~}{𝒪}_o^{}(N_\gamma ,\mathrm{\Lambda },\{r_{kq}\})`$ of elements annihilated by $`W_{kq}^{}W_{kq}.`$ It is clear from the above lemma that this space is stable under $`W_{kq}\mathrm{\Delta }_{r_{kq}^1}W_{kq}.`$ We may interpret what was proved in the lemma as giving a semi-classical normal form for $`\mathrm{\Delta }`$ or $`\mathrm{\Delta }_{r_{kq}^1}`$ in the following sense:
(2.49) TheoremThe $`r_{kq}^1`$-pseudodifferential operator
$$W_{kq}:𝒪_o^{}(N_\gamma ,\mathrm{\Lambda },\{r_{kq}\})\stackrel{~}{𝒪}_o^{}(N_\gamma ,\mathrm{\Lambda },\{r_{kq}\})$$
has the properties:
$$W_{kq}^{}W_{kq}^1IW_{kq}W_{kq}^{}I0$$
$$W_{kq}^{}T_{r_{kq}^1}^1\mathrm{\Delta }T_{r_{kq}^1}W_{kq}^2+f_o(I_{\gamma 1},\mathrm{},I_{\gamma n})+\frac{f_1(I_{\gamma 1},\mathrm{},I_{\gamma n})}{}+\mathrm{}.$$
$$W_{kq}e^{ir_{kq}s}U_q(s,y)=e^{ir_{kq}s}W_{kq}U_q(s,y)\mathrm{\Phi }_{kq}(s,y)$$
(see (2.1)).
## 3. Normal form: Proof of Theorem B
The intertwining operators $`W_{kq}`$ will now be assembled into the Fourier-Hermite -series integral operator
$`(3.1)`$
$$W_\gamma :L^2(S_L^1\times ^n,dsdy)L^2(S_L^1\times ^n,dsdy)$$
$$W_\gamma \underset{(k,q)𝐍^{n+1}}{}\widehat{f}(k,q)e^{ir_{kq}s}U_q(s,y)=\underset{(k,q)𝐍^{n+1}}{}\widehat{f}(k,q)e^{ir_{kq}s}W_{kq}U_q(s,y).$$
Also, the dilation operators will be assembled into the operator
$`(3.2)`$
$$T:L^2(S_L^1\times ^n,dsdy)L^2(S_L^1\times ^n,dsdy)$$
$$T\underset{(k,q)𝐍^{n+1}}{}\widehat{f}(k,q)e^{ir_{kq}s}U_q(s,y)=\underset{(k,q)𝐍^{n+1}}{}\widehat{f}(k,q)e^{ir_{kq}s}U_q(s,\sqrt{r_{kq}}y).$$
By theorem (2.49) we then have, at least formally,
$`(3.3)`$
$$W_\gamma ^1T^1\mathrm{\Delta }TW_\gamma ^2+f_o(I_{\gamma 1},\mathrm{},I_{\gamma n})+\frac{f_1(I_{\gamma 1},\mathrm{},I_{\gamma n})}{}+\mathrm{}.$$
The purpose of this section is to make this equivalence precise. Since it is independent of the model we will carry out the proof in the basic model and use the weightless Fermi normal coordinates $`(s,x)`$. For the sake of notational simplicity the $`\stackrel{~}{W}^{}s`$ and $`\stackrel{~}{Q}`$’s, transferred back to the model in this way, will be written as $`W`$’s and $`Q`$’s. Also in place of the $`r_{kq}^1`$’s we will use the weightless $`h_{kq}`$’s with $`h_{kq}^1=(2\pi k+_{j=1}^n(q_j+\frac{1}{2})\alpha _j)=Lr_{kq}.`$
(3.4) Propostion $`TW_\gamma T^1`$ is a (standard) Fourier integral operator, well-defined and invertible on the microlocal neighborhood (0.1) in $`T^{}(S_L^1\times ^n)`$.
Sketch of Proof:
To simplify, we first consider the unitarily equivalent operator $`\stackrel{~}{T}\stackrel{~}{W}\stackrel{~}{T}^1`$ in the microlocal neighborhood (0.1) in the twisted model, with
$`(3.5)`$
$$\stackrel{~}{W}:_\alpha _\alpha $$
$$\stackrel{~}{W}(e^{i\frac{s}{Lh_{kq}}}\gamma _q):=e^{i\frac{s}{Lh_{kq}}}W_{h_{kq}}\gamma _q,$$
and with $`\stackrel{~}{T}`$ the dilation operator like (3.2) relative to the basis $`e^{i\frac{s}{Lh_{kq}}}\gamma _q`$. We then factor $`\stackrel{~}{T}\stackrel{~}{W}\stackrel{~}{T}^1`$ as the product $`\stackrel{~}{T}\stackrel{~}{W}\stackrel{~}{T}^1=j^{}\stackrel{~}{T}V\stackrel{~}{T}^1`$ where:
$`(3.6)`$
$$V:_\alpha L_{loc}^2(\times \times ^n)$$
$$V=\mathrm{\Pi }_{j=o}^{\mathrm{}}exp[iD_s^{\frac{j}{2}}Q_{\frac{j}{2}}(s^{},y,D_y)]$$
that is,
$$Ve^{i\frac{s}{Lh_{kq}}}\gamma _q(x):=e^{i\frac{s}{Lh_{kq}}}W_{h_{kq}}(s^{},x,D_x)\gamma _q(x),$$
and where
$`(3.7)`$
$$j^{}:C^{\mathrm{}}(\times \times ^n)C^{\mathrm{}}(\times ^n)$$
$$j^{}f(s,x)=f(s,s,x)$$
is the pullback under the partial diagonal embedding. We note that the variable $`s^{}`$ in $`V`$ occurs essentially as a parameter, so we can (and often will) regard V as a one parameter family of operators $`V_s^{}`$ on $`_\alpha .`$
It is not clear that the infinite product in (3.6) is well-defined, nor what kind of operators the factors are. On the second point, we note that $`D_s`$ is elliptic in the set (0.1) in $`T^{}(S_L^1\times ^n)`$ and that, as an operator-valued function of $`s^{}`$,
$`(3.8)`$
$$\stackrel{~}{T}D_s^{\frac{j}{2}}Q_{\frac{j}{2}}(s^{},x,D_x)\stackrel{~}{T}^1C^{\mathrm{}}(S_L^1,\mathrm{\Psi }^1(S_L^1\times ^n)).$$
Here of course $`C^{\mathrm{}}(S_L^1,𝒜)`$ denotes the smooth functions of $`s^{}`$ with values in the algebra $`𝒜`$; to simplify the notation we do not indicate the possible twisting. Indeed, we first observe that conjugation by $`\stackrel{~}{T}`$ converts the mixed polyhomogeneous- isotropic algebra $`\mathrm{\Psi }^{}(S_L^1)𝒲^{}`$ into the pure polyhomogeneous algebra $`\mathrm{\Psi }^{}(S_L^1)\mathrm{\Psi }^{}(^n)`$. This essentially follows from the fact that
$$\stackrel{~}{T}a^w(x,D_x)\stackrel{~}{T}^1=a^w(D_s^{\frac{1}{2}}x,D_s^{\frac{1}{2}}D_x)$$
(see \[G.2\] or \[CV\] for further details). Obviously, $`\stackrel{~}{T}D_s\stackrel{~}{T}^1\mathrm{\Psi }^1(S_L^1\times ^n)`$, and by (2.44(vi)) we also have
$$\stackrel{~}{T}Q_{\frac{j}{2}}\stackrel{~}{T}^{}C^{\mathrm{}}(S_L^1,\mathrm{\Psi }^{j+1}(S^1\times ^n)).$$
A simple calculation shows further that $`\stackrel{~}{T}D_s^{\frac{j}{2}}Q_{\frac{j}{2}}(s^{},x,D_x)\stackrel{~}{T}^1`$ is a one-parameter family of pseudodifferential operators of real principal type over $`S_L^1\times ^n`$ whose principal symbols have nowhere radial Hamilton vector field. It follows in a well-known way that these operators are microlocally conjugate to $`D_1`$ (see \[Ho IV, Proposition 26.1.2\]) and hence that their exponentials give a smooth one-parameter family of Fourier Integral operators. Consequently, after conjugation with $`\stackrel{~}{T}`$, the factors in (3.6) are for each $`s^{}S^1`$ microlocally well-defined Fourier Integral operators in the neighborhood (0.1), with smooth dependence on $`s^{}`$. As for the infinite product, we note that the principal symbol of $`\stackrel{~}{T}D_s^{\frac{j}{2}}Q_{\frac{j}{2}}\stackrel{~}{T}^1`$ has the form $`\sigma ^{\frac{j}{2}}q_{\frac{j}{2}}(s^{},\sqrt{\sigma }y,\frac{1}{\sqrt{\sigma }}\eta )`$, where $`q_{\frac{j}{2}}`$ is a homogenous polynomial in $`(y,\eta )`$ of degree j+2 (2.36(iv)). Its exponential can therefore be constructed microlocally in (0.1) as a one parameter family of Fourier integral operators over $`S_L^1\times ^n`$ with phase functions of the form
$$\mathrm{\Psi }_{\frac{j}{2}}(s^{},s,s^{\prime \prime }\sigma ,x,x^{\prime \prime },\xi )=(ss^{\prime \prime })\sigma +xx^{\prime \prime },\xi +\psi _{\frac{j}{2}}(s^{},s,s^{\prime \prime },\sigma ,x,x^{\prime \prime },\xi )$$
with
$`(3.9b)`$
$$\psi _{\frac{j}{2}}O_{j+2}S_{cl}^1$$
Here, $`S_{cl}^1`$ denotes the space of classical symbols of first order in (0.1) and $`O_k`$ denotes the elements which vanish to order k along $`(x,\xi )=(0,0).`$ (Henceforth we will use the symbol $`O_k`$ for objects of any kind which vanish to order k on this set.) It follows that for fixed $`s^{}`$ and sufficiently small $`ϵ`$, the phase parametrizes the graph of a homogeneous canonical transformation of $`T^{}(S_L^1\times ^n)`$ which is $`C^{j+1}`$\- close to the identity. Any finite product $`\mathrm{\Pi }_{j=o}^N`$ in (3.7) is therefore a clean composition and the phase $`\mathrm{\Psi }_N`$ of the product parametrizes the corresponding compostion $`\chi _N`$ of canonical transformations. By (3.9b), we have $`\chi _{N+1}\chi _NO_{N+1}`$ and so there exists a smooth one parameter family homogeneous canonical transformations $`\chi _s^{}`$, equal to the identity along $`(x,\xi )=(0,0)`$ such that $`\chi _s^{}\chi _{s^{},N}O_{N+1}`$ and with a generating function $`\psi _{\mathrm{}}`$ satisfying $`\psi _{\mathrm{}}\psi _NO_N`$ (cf. \[Sj, Proposition 1.3\]). Regarding the amplitudes, the situation is of a similar kind. Denoting the amplitude of $`\stackrel{~}{T}exp(D_s^{\frac{j}{2}}Q_{\frac{j}{2}})\stackrel{~}{T}^1`$ by $`a_{\frac{j}{2}}(s^{},s,\sigma ,y\eta )`$ we have
$`(3.10)`$
$$a_{\frac{j}{2}}1modO_{j+2}S^o.$$
It follows that the amplitudes $`a_N`$ of the products $`\mathrm{\Pi }_{j=o}^N`$ satisfy $`a_{N+1}a_NO_{N+1}S^o`$ and hence there exists an element $`a_{\mathrm{}}`$ satisfying $`a_{\mathrm{}}a_NO_{N+1}S^o`$. The pair $`(a_{\mathrm{}},\psi _{\mathrm{}})`$ then determine a one-parameter family of Fourier integral operators of order zero which agrees with $`\mathrm{\Pi }_{j=o}^N`$ up to an error which also vanishes to order $`N`$. Finally, composing with $`j^{}`$ just sets $`s^{}=s`$ in the kernel, which visibly remains Fourier integral, with phase function parametrizing the graph of a canonical transformation $`\chi .`$ Invertibility then follows from the fact that $`\stackrel{~}{W}^{}W`$ is a pseudodifferential operator with principal symbol the square of $`j^{}(\sigma _V)`$, which is easily seen to equal 1.
The proposition follows by expressing
$$TW_\gamma T^1=T\mu (a)^{}\stackrel{~}{T}^1\stackrel{~}{T}W\stackrel{~}{T}^1\stackrel{~}{T}\mu (a)T^1$$
and noting that $`\stackrel{~}{T}\mu (a)T^1`$ is also a standard Fourier Integral operator. ∎
We now give the proof of the quantum normal form Theorem B for $`\sqrt{\mathrm{\Delta }}`$, stated in an equivalent form in terms of $`W_\gamma .`$ As in the introduction, the notation $`AB`$ means that the complete (Weyl) symbol of $`AB`$ vanishes to infinite order at $`\gamma .`$ We also use the notation $`O_j\mathrm{\Psi }^m`$ for pseudodifferential operators of order m whose Weyl symbols vanish to order j at $`(y,\eta )=(0,0)`$. Here, pseudodifferential operator could mean of the standard polyhomogeneous kind, or of the mixed polyhomogeneous-isotropic kind as in $`\mathrm{\Psi }^k(S_L^1)𝒲^l`$, in which case the total order is defined to be $`m=k+l`$. To simplify notation, we will denote the space of mixed operators of order m by $`\mathrm{\Psi }_{mx}^m(S_L^1\times ^n)`$.
(3.11) Theorem B Let $`TW_\gamma T^1`$ be the Fourier Integral operator of Proposition (3.4), defined over a conic neighborhood of $`R^+\gamma `$ in $`T^{}(S_L^1\times ^n)`$. Then:
$$W_\gamma ^1T^1\sqrt{\mathrm{\Delta }}TW_\gamma P_1(,I_{\gamma 1},\mathrm{},I_{\gamma n})+P_o(,I_{\gamma 1},\mathrm{},I_{\gamma n})+\mathrm{},$$
where
$`(3.14b)`$
$$P_1(,I_{\gamma 1},\mathrm{},I_{\gamma n})+\frac{p_1^{[2]}(I_{\gamma 1},\mathrm{},I_{\gamma n})}{L}+\frac{p_2^{[3]}(I_{\gamma 1},\mathrm{},I_{\gamma n})}{(L)^2}+\mathrm{}$$
$$P_m(,I_{\gamma 1},\mathrm{},I_{\gamma n})\underset{k=m}{\overset{\mathrm{}}{}}\frac{p_k^{[km]}(I_{\gamma 1},\mathrm{},I_{\gamma n})}{(L)^j}$$
with $`p_k^{[km]}`$, for m=-1,0,1,…, homogenous of degree l-m in the variables $`(I_{\gamma 1},\mathrm{},I_{\gamma n})`$ and of weight -1.
Proof:
As a semi-classical expansion in the “parameter” $`h=\frac{1}{L}`$, (3.3) may be rewritten in terms of $`\sqrt{\mathrm{\Delta }}`$ :
$`(3.12)`$
$$W_\gamma ^1T^1\sqrt{\mathrm{\Delta }}TW_\gamma +\frac{p_1(I_{\gamma 1},\mathrm{},I_{\gamma n})}{L}+\frac{p_2(I_{\gamma 1},\mathrm{},I_{\gamma n})}{(L)^2}+\mathrm{}.$$
From the fact that the numerators $`f_j(I_{\gamma 1},\mathrm{},I_{\gamma n})`$ in (3.3) are polynomials of degree j+2 (2.36viii) and of weight -2, the numerators $`p_k(I_{\gamma 1},\mathrm{},I_{\gamma n})`$ are easily seen to be polynomials of degree $`k+1`$ and of weight -1. Hence they may be expanded in homogeneous terms
$`(3.13)`$
$$p_k=p_k^{[k+1]}+p_k^{[k]}+\mathrm{}p_k^{[o]},$$
with $`p_k^{[j]}`$ the term of degree j and still of weight -1. The right side of (3.12) can then be expressed as a sum of homogeneous operators:
$`(3.14a)`$
$$P_1(,I_{\gamma 1},\mathrm{},I_{\gamma n})+P_o(,I_{\gamma 1},\mathrm{},I_{\gamma n})+\mathrm{}$$
with
$`(3.14b)`$
$$P_1(,I_{\gamma 1},\mathrm{},I_{\gamma n})+\frac{p_1^{[2]}(I_{\gamma 1},\mathrm{},I_{\gamma n})}{L}+\frac{p_2^{[3]}(I_{\gamma 1},\mathrm{},I_{\gamma n})}{(L)^2}+\mathrm{}$$
$$P_m(,I_{\gamma 1},\mathrm{},I_{\gamma n})\underset{k=m}{\overset{\mathrm{}}{}}\frac{p_k^{[km]}(I_{\gamma 1},\mathrm{},I_{\gamma n})}{(L)^k}.$$
The remainder term in (3.12) can be described as follows:
$`(3.15)`$
$$W_\gamma ^1T^1\sqrt{\mathrm{\Delta }}TW_\gamma [+\frac{p_1(I_{\gamma 1},\mathrm{},I_{\gamma n})}{L}+\frac{p_2(I_{\gamma 1},\mathrm{},I_{\gamma n})}{(L)^2}+\mathrm{}+\frac{p_m(I_{\gamma 1},\mathrm{},I_{\gamma n})}{(L)^m}]$$
$$_{k=o}^{m+1}O_{2(m+1k)}\mathrm{\Psi }_{mx}^{1k}(S_L^1\times ^n).$$
The remainder terms in (3.14b) are given by:
$`(3.16a)`$
$$P_1(,I_{\gamma 1},\mathrm{},I_{\gamma n})[+\frac{p_1^{[2]}(I_{\gamma 1},\mathrm{},I_{\gamma n})}{L}+\mathrm{}+\frac{p_N^{[N+1]}(I_{\gamma 1},\mathrm{},I_{\gamma n})}{(L)^k}]$$
$$O_{2(N+2)}\mathrm{\Psi }_{mx}^1(S_L^1\times ^n)$$
$`(3.16b)`$
$$P_m(,I_{\gamma 1},\mathrm{},I_{\gamma n})\underset{k=m}{\overset{N}{}}\frac{p_k^{[km]}(I_{\gamma 1},\mathrm{},I_{\gamma n})}{(L)^k}O_{2(N+1m)}\mathrm{\Psi }_{mx}^m(S_L^1\times ^n).$$
Hence the expansion (3.12) is also asymptotic in the sense of $`.`$ For the statement of Theorem B in the introduction, we only need to conjugate under $`\mu (\stackrel{~}{a}).`$
Remark The remainder in (3.15) could be equivalently described by saying that its complete symbol lies in the symbol class $`S_{cl}^1S^{1,2(m+1)}`$ of Boutet-de-Monvel \[BM\].
## 4. Residues and Wave Invariants: Proof of Theorem C
The aim now is to use the normal form of $`\sqrt{\mathrm{\Delta }}`$ near $`\gamma `$ to calculate the wave invariants $`a_{\gamma k}`$ associated to $`\gamma `$. In terms of the model (see the statement of Theorem B in the introduction) we may write the normal form as
$`(4.1)`$
$$𝒟=\overline{\psi }(I_1,\mathrm{},I_n)[𝒟_N+_N]$$
with
$$𝒟_N:=+\frac{p_1(I_1,\mathrm{},I_n)}{L}+\mathrm{}+\frac{p_N(I_1,\mathrm{},I_n)}{(L)^N}$$
and with
$$_N_{k=o}^{N+1}O_{2(N+1k)}\mathrm{\Psi }_{mx}^{1k}(S_L^1\times ^n).$$
Our first observation is that we can drop a sufficiently high remainder term $`_N`$ in the calculation of a given wave invariant. We prove this by working out, roughly, which parts of the complete symbol of a general first order pseudodifferential operator $`P`$ of real principal type contribute to the wave invariant of a given order associated to a non-degenerate closed bicharacteristic $`\gamma .`$
In the following proposition, $`(s,y)`$ will denote any coordinates in a tubular neighborhood $`U`$ of the projection of $`\gamma `$ to M with the property that the equation for $`^+\gamma `$ in the associated symplectic coordinates $`(s,\sigma ,y,\eta )`$ is $`y=\eta =0.`$ Also, $`Op`$ we will denote a fixed quantization of symbols in a conic neighborhood (0.1) of $`^+\gamma `$ to pseudodifferential operators, and $`p(s,\sigma ,y,\eta )p_1+p_o+\mathrm{}`$ will denote the complete symbol of $`P`$. The Taylor expansion of $`p_j`$ at $`^+\gamma `$ will be written in the form,
$$p_j(s,\sigma ,y,\eta )=\sigma ^jp_j(s,1,y,\frac{\eta }{\sigma })=\sigma ^j(p_j^{[o]}+p_j^{[1]}+\mathrm{})$$
with $`p_j^{[m]}(s,1,y,\frac{\eta }{\sigma })`$ homogeneous of degree $`m`$ in $`(y,\frac{\eta }{\sigma }).`$ We set $`P_j:=Op(p_j),P_j^{[m]}:=Op(p_j^{[m]})`$ and $`P_j^N=_{mN}P_j^{[m]}`$. Finally, we will denote by $`\tau _{\gamma k}(P)`$ the coefficient of $`(tL+i0)^klog(tL+i0)`$ (or of $`(tL+i0)^1`$ in the case of $`k=1`$) in the singularity expansion of the microlocal unitary group $`e^{itP}`$ near $`\gamma `$.
(4.2) Proposition
$$\tau _{\gamma k}(P)=\tau _{\gamma k}(P_1^{2(k+2)}+P_o^{2(k+1)}+\mathrm{}P_{k1}^o).$$
Proof: As mentioned several times, the wave invariant $`\tau _{\gamma k}(P)`$ is given by the the non-commutative residue
$`(4.3)`$
$$\tau _{\gamma k}(P)=resD_t^ke^{itP}|_{t=L}=resP^ke^{itP}|_{t=L}$$
(see \[Z.1\] or \[G.2\]). We first describe how (4.3) leads to a local formula for $`\tau _{\gamma k}(P)`$ in terms of the jets of the amplitudes and phases of a microlocal parametrix
$`(4.4)`$
$$F_\gamma (t,x,x^{})=_^ne^{i\varphi (t,x,x^{},\theta )}a(t,x,x^{},\theta )\psi (x,\theta )𝑑\theta $$
for $`e^{itP}`$ near $`\{L\}\times R^+\gamma \times R^+\gamma `$. The remainder of the proof will connect this data to the complete symbol of $`P`$.
The amplitude $`a`$ in (4.4) is a classical symbol, hence has the expansion $`a_{j=0}^{\mathrm{}}a_j`$ with $`a_j`$ homogeneous of degree $`j`$ for $`|\theta |1`$. Since the residue or wave invariant depends only on the wave kernel in a microlocal neighborhood of $`\{L\}\times \gamma \times \gamma `$ it is unchanged by the insertion homogeneous cut-off $`\psi `$, supported in (0.1) and equal to one on a smaller conic neighborhood of $`\gamma `$. The phase $`\varphi (t,x,x^{},\theta )`$ parametrizes a piece ofthe graph of the Hamilton flow $`exptH_{p_1}`$ of $`p_1`$ near $`\{L\}\times \gamma ,`$ and may be assumed to consist of only $`n`$ phase variables. Writing $`x=(s,y)`$ and $`\theta =(\sigma ,\eta )`$ as above and choosing $`\psi `$ of the form $`\psi (\frac{\eta }{\sigma })`$, we get
$`(4.5)`$
$$\tau _{\gamma k}(P)=Res_{s=0}_^+_{^{n1}}_Ua^{(k)}(L,s,y,s,y,\sigma ,\eta )\psi (\frac{\eta }{\sigma })e^{i\varphi (L,s,y,s,y,\sigma ,\eta )}|\sigma |^s𝑑\sigma 𝑑\eta 𝑑v(s,y)$$
where for certain universal constants $`C_{\alpha \beta }`$,
$`(\mathrm{4.5.1})`$
$$a^{(k)}:=\underset{|\alpha |+|\beta |=k}{}C_{\alpha \beta }(D_t^{\alpha _1}\varphi )^{\beta _1}\mathrm{}(D_t^{\alpha _n}\varphi )^{\beta _n}D_t^{\alpha _{n+1}}a.$$
We note here that the residue can be calculated using any gauging of the trace, in particular by powers $`|\sigma |^s`$ of the elliptic symbol $`\sigma `$ in (0.1). Changing variables to $`\stackrel{~}{\eta }:=\frac{\eta }{\sigma }`$ and denoting by $`a_{kj}^{(k)}`$ the term of order $`kj`$ in the polyhomogenous expansion of the kth order amplitude $`a^{(k)}`$, and by $`B_ϵ`$ the ball of radius $`ϵ`$ in $`_{\stackrel{~}{\eta }}^n`$, we get
$`(4.6)`$
$$\tau _{\gamma k}(P)=$$
$$\underset{j=o}{\overset{\mathrm{}}{}}Res_{s=0}_U_{R^+}_{B_ϵ}e^{i\sigma \varphi (L,s,y,s,y,1,\stackrel{~}{\eta })}a_{kj}^{(k)}(L,s,y,s,y,1,\stackrel{~}{\eta })\psi (s,y,\stackrel{~}{\eta })\sigma ^{n+k1js}𝑑\sigma 𝑑\stackrel{~}{\eta }𝑑v(s,y)$$
$$=\underset{j=0}{\overset{\mathrm{}}{}}_U_{B_ϵ}\psi a_j^{(k)}(\varphi (L,s,y,s,y,1,\stackrel{~}{\eta })+i0)^{s+jnk}𝑑\stackrel{~}{\eta }𝑑v(s,y).$$
The non-degeneracy of $`\gamma `$ as a closed geodesic implies its non-degeneracy as the critical manifold of $`\varphi (s,y,s,y,1,\stackrel{~}{\eta })`$. Hence we have
$$\varphi (L,s,y,s,y,1,\stackrel{~}{\eta })=\varphi (L,s,0,s,0,1,0)+Q_s(y,\stackrel{~}{\eta })+g(s,y,\stackrel{~}{\eta })$$
with $`Q`$ a non-dengenerate symmetric bilinear form and with $`g=0_s(|(y,\stackrel{~}{\eta })|^3).`$ Expanding
$$(Q_s(y,\stackrel{~}{\eta })+i0+g)^\lambda =\underset{p=o}{\overset{\mathrm{}}{}}C_{(\lambda ,p)}(Q_s(y,\stackrel{~}{\eta })+i0)^{\lambda p}g^p$$
we get that
$`(4.7)`$
$$\tau _{\gamma k}(P)=\underset{j,p=o}{\overset{\mathrm{}}{}}C_{(s+jnk,p)}Res_{s=0}_\gamma _{|y|<ϵ}_{B_ϵ}(Q_s(y,\stackrel{~}{\eta })+i0)^{s+jnkp}(g^pa_{kj}^{(k)}\psi J(s,y))𝑑s𝑑y𝑑\stackrel{~}{\eta }$$
with $`J(s,y)`$ the volume density and with binomial coefficients $`C_{(\lambda ,p)}`$. The family of distributions$`(Q_s(y,\stackrel{~}{\eta })+i0)^\lambda `$ is meromorphic with simple poles at the points $`\lambda =(n1)r,r=0,1,2,\mathrm{}`$ and with residue $`C_r(Q_s^1(D_y,D_\eta ))^r\delta (y,\eta )`$ for certain constants $`C_r`$ (see, \[G.S., p.276\]). Here we have dropped the tilde in the notation for $`\eta .`$ For $`s=0`$ there are possible poles when $`k+p+1j0`$ with residues
$`(4.8)`$
$$C_rC_{(jnk,p)}_\gamma Q_s^1(D_y,D_\eta )^{k+p+1j}(g^pa_{kj}^{(k)}J)|_{(y=\eta =0)}ds.$$
However, the residue vanishes unless $`2(k+p+1j)3p`$ since $`g^p`$ vanishes to order $`3p,`$ constraining the sum in (4.6) to $`2j+p2(k+1).`$ Hence we have, with new constants $`C_{jkp}`$,
$`(4.9)`$
$$\tau _{\gamma k}(P)=\underset{j,p:2j+p2(k+1)}{}C_{jkp}_\gamma Q_s^1(D_y,D_\eta )^{k+p+1j}(g^pa_{kj}^{(k)}J)|_{(y,\eta )=(0,0)}ds.$$
It follows first that $`jk+1`$, hence that only the terms $`a_k^{(k)},\mathrm{},a_1^{(k)}`$ in the amplitude $`a^{(k)}`$ contribute to (4.9). To determine the corresponding range of terms $`a_j`$ in the amplitude $`a`$, we observe from (4.5.1) that
$`(4.10)`$
$$a_{kj}^{(k)}:=\underset{|\alpha |+|\beta |=k}{}C_{\alpha \beta }(D_t^{\alpha _1}\varphi )^{\beta _1}\mathrm{}(D_t^{\alpha _n}\varphi )^{\beta _n}D_t^{\alpha _{n+1}}a_{kj|\beta |},$$
hence that only the terms $`a_o,a_1,\mathrm{},a_{k1}`$ contribute to (4.9).
The same kind of calculation determines $`resAe^{itP}|_{t=L}`$, for any $`k`$th order pseudodifferential operator $`A`$ supported microlocally in (0.1) in terms of the complete symbol of $`A`$ and the phases and amplitudes in (4.3). In particular, we see that $`resD_t^kAe^{itP}|_{t=L}=0`$ if $`A`$ has order $`(k+2)`$. This implies:
$`(4.11)`$
$$resD_t^ke^{itP}=resD_t^ke^{it(P_1+\mathrm{}P_{(k+1)})}$$
since $`e^{itP}=A_k(t)e^{it(P_1+\mathrm{}+P_{(k+1)})}`$ with $`A_k(t)I\mathrm{\Psi }^{k2}.`$ Indeed, with $`H_k=P_1+\mathrm{}P_{(k+1)},V_k=PH_k,\stackrel{~}{B}(t):=e^{itH_k}B(t)e^{itH_k}`$, we have
$$D_t\stackrel{~}{A}_k=\stackrel{~}{V}_k(t)\stackrel{~}{A}_k(t),\stackrel{~}{A}_k(0)=I.$$
Thus $`\stackrel{~}{A}_k(t)`$ is the time-ordered exponential of a pseudodifferential operator of order $`k2`$, hence equal to $`I`$ modulo $`\mathrm{\Psi }^{k2}`$; consequently so is $`A_k(t).`$
We next consider which jets $`j_\gamma ^ma_i`$ of $`a_o,\mathrm{}a_{k1}`$ and $`j_\gamma ^m\varphi `$ of $`\varphi `$ contribute to (4.9). Here by $`j_\gamma ^ma_i`$ we mean the Taylor polynomial of degree $`m`$ of $`a_i(t,s,y,s,y,\sigma ,\eta )`$ in the $`(y,\eta )`$ variables for $`t`$ near $`L`$. From (4.9-10) it is evident that the maximum number of $`(y,\eta )`$ derivatives on $`a_i`$ occurs when $`p=0`$ in (4.9) and when $`|\beta |=k,kj|\beta |=i`$ in (4.10). It follows that at most $`2(k+1i)`$ derivatives fall on $`a_i`$ in (4.9). Also, the maximum number of $`(y,\eta )`$ derivatives on the phase occurs in terms where a factor of $`g`$ is differentiated the maximum number of times; namely in terms with $`p=1`$ and with $`j=0,`$ in which the phase is differentiated $`2(k+2)`$ times. Hence
$`(4.12)`$
$$\tau _{\gamma k}(P)\text{depends only on}j_\gamma ^{2(k+1)}a_o,j_\gamma ^{2k}a_1,\mathrm{}j_\gamma ^oa_{(k+1)};j_\gamma ^{2(k+2)}\varphi .$$
To give bounds on the jets of the terms $`p_j(j=1,0,\mathrm{}(k+1))`$ in the symbol of $`P`$ which contribute to the jets $`j_\gamma ^{2(k+1i)}a_i`$, we now have to consider some details of the construction of the parametrix (4.3). We first recall that the amplitudes are obtained by solving transport equations of the form
$`(4.13)`$
$$D_ta_j=\underset{m=0}{\overset{j}{}}[𝒫_{1,m+1}a_{j+m}+\underset{\nu =o}{\overset{jm}{}}𝒫_{\nu ,m}a_{j+m+\nu }]$$
where $`𝒫_{\nu ,m}=𝒫_{\nu ,m}(\varphi ,s,y,D_s,D_y)`$ is a differential operator of order $`m`$ obtained from $`Op(p_\nu )`$ as follows:
$`(\mathrm{4.14.1})`$
$$e^{i\rho \varphi }Op(p_\nu )e^{i\rho \varphi }\underset{m=0}{\overset{\mathrm{}}{}}\rho ^{\nu m}𝒫_{\nu ,m}$$
$`\varphi `$ being the phase in the microlocal parametrix (see \[Tr, Ch.VI,(5.28)\]). Consistently with (4.11), only the terms $`p_1,p_o,\mathrm{},p_{(k+1)}`$ in the symbol of $`P`$ contribute to the transport equation for $`a_i(i=o,\mathrm{},(k+1).`$ To determine which jet of $`p_\nu `$ is involved in $`𝒫_{\nu ,m}`$, we further recall \[loc.cit.,Theorem 3.1\] that the expansion in (4.14.1) is obtained by re-arranging terms in the expansion
$`(\mathrm{4.14.2})`$
$$e^{i\rho \varphi }Op(p_\nu )e^{i\rho \varphi }u\underset{\alpha }{}\frac{1}{\alpha !}_\xi ^\alpha p_\nu (s,y,\rho d\varphi )𝒩_\alpha (\varphi ;\rho ,D_s,D_y)u$$
with
$$𝒩_\alpha (\varphi ;\rho ,D_x)u=D_x^{}^\alpha e^{i\varphi _2(x,x^{})}u(x^{})|_{x=x^{}}$$
$$\varphi _2(x,x^{})=\varphi (x)\varphi (x^{})xx^{},\varphi (x).$$
In view of the fact that $`𝒩_\alpha `$ is a polynomial in $`\rho `$ of degree $`\frac{|\alpha |}{2}`$, we see that $`|\alpha |2m`$ in any term of (14.4.2) contributing to $`𝒫_{\nu ,m}`$. Hence $`𝒫_{\nu ,m}`$ involves at most $`2m`$ derivatives of $`p_\nu `$, and so the transport equation for $`a_i`$ involves at most $`2(i\nu )`$ derivatives of $`p_\nu `$ ($`\nu =1,0,\mathrm{},i.`$)
To draw these kinds of conclusions about the amplitudes $`a_i`$ themselves, as opposed to the coefficients in their transport equations, we have to take into account the initial conditions in the transport equations. Unfortunately, the defining initial conditions occur at $`t=0`$ while we are interested in the long time behaviour at $`t=L`$. Were there no conjugate pairs along geodesics near $`\gamma `$, the transport equations could be solved up to $`t=L`$ and since the initial conditions can be taken in the form $`a_o|_{t=0}1,a_i|_{t=0}0`$ we could conclude that the $`a_i`$’s are integrals of polynomials in at most $`2(i\nu )`$ derivatives of $`p_\nu `$ and in at most $`2(i+1)`$ derivatives of the phase (coming from the $`𝒩_\alpha `$’s). However, in the case of elliptic closed geodesics, there will always be conjugate pairs for $`L`$ sufficiently large and we cannot construct the parametrix so simply. Rather we will use the group property $`U(L)=U(L/N)^N`$ with N sufficiently large that a parametrix for $`U(L/N)`$ can be constructed by the geometric optics method. We then have to eliminate all but $`n`$ phase variables in the power, which will lead to crude but serviceable bounds on the order of the jets.
We therefore begin with the construction of a short time parametrix $`F_{\delta k}`$, valid for $`|t|<\delta ,`$ of the form (4.4) with phase $`\varphi (t,x,x^{},\theta )=S(t,x,\theta )x^{},\theta ,`$ satisfying $`_tS+p_1(x,d_xS)=0,S|_{t=o}=x,\theta `$, and with amplitude $`a_{j=0}^{k+1}a_j`$ satisfying (4.13) on $`|t|\delta `$ and with initial conditions $`a_o|_{t=0}=1,a_j|_{t=0}=0(j>0)`$. By the observations above, we have that $`a_{\delta ,i}(t,x,x^{},\theta )`$ involves at most $`2(i\nu )`$ derivatives of $`p_\nu (\nu i)`$ for $`t`$ in this interval. As is well-known, we then have $`e^{itP}F_{\delta k}(t)=R_{ok}(t)`$ with the remainder $`R_{\delta k}(t)`$ a Fourier integral operator of order $`(k+2)`$ associated to the graph of $`exptH_{p_1}`$ for $`|t|<\delta `$. Since $`F_{\delta k}`$ is of order zero, it follows that for any $`N`$ and for $`|t|<\delta ,`$ $`e^{itNP}F_{\delta k}^N(t)`$ modulo Fourier integral operators of order $`k2`$ in this class. Hence, the desired parametrix (4.4) can be constructed by choosing $`N`$ so that $`L/N<\delta `$ and by eliminating all but $`n`$ phase variables in $`F_{\delta k}^N`$ by the stationary phase method. Just as in (4.11) the remainder terms with a factor of $`R_{\delta k}`$ will not contribute to $`\tau _{\gamma k}(P)`$, nor will terms $`a_i`$ in the final amplitude with $`i>k+1.`$
To complete the proof of the proposition, we have to count the number of $`(y,\eta )`$-derivatives of $`p_\nu `$ which appear in the terms
$`(\mathrm{4.15.1})`$
$$D_t^{k|\beta |}Q_s^1(D_y,D_\eta )^{k+1j}\widehat{a}_{kj|\beta |}|_{(y,\eta )=(0,0)},$$
and the number of derivatives of the final phase (hence of $`p_1`$) that occur in the terms
$`(\mathrm{4.15.2})`$
$$Q_s^1(D_y,D_\eta )^{k+p+1j}g^p(D_t^{\alpha _1}\varphi )^{\beta _1}\mathrm{}(D_t^{\alpha _n}\varphi )^{\beta _n}D_t^{\alpha _{n+1}}\widehat{a}_{kj|\beta |}|_{(y,\eta )=(0,0)}$$
with $`\widehat{a}_i`$ the terms in the final amplitude. We start from the fact that, in an obvious notation,
$`(4.16)`$
$$F_{k\delta }^N=_{^{nN}}_{M^{N1}}e^{i(\varphi _1+\mathrm{}\varphi _N)}\underset{(i_1,\mathrm{},i_N)𝐍^N}{}a_{i_1}\mathrm{}a_{i_N}d\theta _1\mathrm{}d\theta _Ndv(x_1)\mathrm{}dv(x_{N1})$$
and that the elimination of all but $`n`$ phase variables by the stationary phase method gives essentially the same formulae for the final amplitudes $`\widehat{a}_i`$ as in the case of non-homogeneous Lagrangean distributions of the form $`u(x,y)e^{i\omega f(x,y)}𝑑x`$ with $`f(x,y)`$ a real-valued function defined near $`(0,0)`$, with $`f_x^{}(0,0)=0`$, and with $`f_{xx}^{\prime \prime }(0,0)`$ non-singular. In this situation, we have (see \[HoI, Theorems 7.7.5-6\])
$$u(x,y)e^{i\omega f(x,y)}𝑑x=C\frac{e^{i\omega f(x(y),y)}}{|det(\omega f_{xx}^{\prime \prime }(x(y),y))|^{\frac{1}{2}}}\underset{j+\frac{n}{2}<M}{}L_{f,j,y}u(x(y),y)\omega ^j$$
modulo terms of order $`0(\omega ^{(M)})`$, with $`L_{f,j,y}`$ a differential operator of order 2j whose coefficients involve at 2j derivatives of $`f^{\prime \prime }`$. It follows that $`\widehat{a}_i`$ has the form (with $`I=(i_1,\mathrm{},i_N),\mathrm{\Phi }=\varphi _1+\mathrm{}\varphi _N)`$)
$$\underset{I,j:|I|+j=i}{}L_{\mathrm{\Phi },j}a_{i_1}\mathrm{}a_{i_N}.$$
Since also $`a_i`$ has the form
$$F_i(p_1,Dp_1,\mathrm{},D^{2(i+1)}p_1,p_o,\mathrm{},D^{2i}p_o,\mathrm{},p_i;\varphi ,\mathrm{},D^{2i}\varphi )$$
with $`F_i`$ a multiple integral of polynomials, and with $`D^k`$ some differential operators of degree k, we see that also $`\widehat{a}_i`$ has the form
$$_i(p_1,\mathrm{},D^{2(i+1)},p_o,\mathrm{},D^{2i}p_o,\mathrm{},D^2p_{i+1},p_i;\varphi ,\mathrm{},D^{2(i+1)}\varphi ).$$
Hence as regards number of derivatives of $`p_1,\mathrm{},p_i,\varphi ,`$ the $`\widehat{a}_i`$’s behave exactly as the $`a_i`$’s, that is, involve at most $`2(i\nu )`$ derivatives of $`p_\nu .`$ Hence (4.9)-(4.10) apply to the final amplitudes, and we conclude that
$$\tau _{\gamma k}(P)\text{depends only on}j_\gamma ^{2(k+2)}p_1,j_\gamma ^{2(k+1)}p_o,\mathrm{}j_\gamma ^op_{(k+1)};j_\gamma ^{2(k+2)}\varphi .$$
Since $`j^{2(k+2)}\varphi `$ depends only on $`j^{2(k+1)}p_1`$, the proof of the proposition is complete..∎
(4.16) Corollary
$$\tau _{\gamma k}(𝒟)=\tau _{\gamma k}(𝒟_{k+1}).$$
The following Lemma combines the previous results in a form applicable to the calculation of the wave invariants. The notation $`TrAP^s`$ (with $`P\mathrm{\Psi }^1`$ elliptic) is short for the zeta function obtained by meromorphic continuation of the trace from $`Res>>0`$ to $`.`$
(4.17) Lemma
$$\tau _{\gamma k}(\sqrt{\mathrm{\Delta }})=Res_{s=0}D_t^kTr\overline{\psi }(,I_1,\mathrm{},I_n)e^{it𝒟_{k+1}}^s|_{t=L}$$
Proof Since $``$ is elliptic in the essential support of $`\overline{\psi }`$, the trace on the right side is well defined and has a meromorphic continuation to $``$ (cf. \[G.2\], \[Z.1,5\]). The residue is of course $`\tau _{\gamma k}(𝒟)`$ by the previous proposition. Hence it suffices to show that
$`(4.18)`$
$$\tau _{\gamma k}(\sqrt{\mathrm{\Delta }})=\tau _{\gamma k}(𝒟).$$
This however follows from the previous proposition combined with Theorem B: Indeed, the proposition shows that $`\tau _{\gamma k}(A)=\tau _{\gamma k}(B)`$ if $`AB`$ in the sense of Theorem B. Also, $`\tau _{\gamma k}(WAW_1)=\tau _{\gamma k}(A)`$ if $`W_1`$ is a parametrix for $`W`$ on the essential support of $`A`$, as follows from the tracial property $`resWAW_1=resW_1WA`$ of the residue (see \[Z.1,5\] for instance).∎
We come now to the calculation of the wave invariants as residue traces of the normal form wave group. But we will simplify (4.17) further before evaluating the residue trace. First, we rewrite everything in terms of $`D_s`$, and $`H_\alpha `$ using (0.3-4). Since
$$\frac{p_\nu (I_1,\mathrm{},I_n)}{(L)^\nu }=\frac{p_\nu (I_1,\mathrm{},I_n)}{(LD_s)^\nu }(I\nu \frac{H_\alpha }{LD_s}+\frac{1}{2}\nu (\nu 1)(\frac{H_\alpha }{LD_s})^2+\mathrm{})$$
and since we can drop the $`D_s^{(k+1)+\nu }H_\alpha ^{k+1\nu }`$ and higher terms by Proposition (4.2), $`𝒟_{k+1}`$ can be written in the form
$`(4.19)`$
$$𝒟_{k+1}D_s+H_\alpha +\frac{\stackrel{~}{p}_1(I_1,\mathrm{},I_n,L)}{LD_s}+\frac{\stackrel{~}{p}_2(I_1,\mathrm{},I_n,L)}{(LD_s)^2}+\mathrm{}+\frac{\stackrel{~}{p}_{k+1}(I_1,\mathrm{},I_n,L)}{(LD_s)^{k+1}}$$
modulo terms which make no contribution to $`\tau _{k\gamma }`$. Secondly, we can use $`LD_s`$ rather than $``$ as the gauging elliptic operator in (4.17). To simplify the notation we will denote all but the first two terms on the right side of (4.19) by
$`(4.20)`$
$$𝒫_{k+1}(D_s,I_1,\mathrm{},I_n,L):=\frac{\stackrel{~}{p}_1(I_1,\mathrm{},I_n,L)}{LD_s}+\frac{\stackrel{~}{p}_2(I_1,\mathrm{},I_n,L)}{(LD_s)^2}+\mathrm{}+\frac{\stackrel{~}{p}_{k+1}(I_1,\mathrm{},I_n,L)}{(LD_s)^{k+1}}.$$
Then we have:
$`(4.21)`$
$$\tau _{\gamma k}(\sqrt{\mathrm{\Delta }})=Res_{z=0}TrD_t^k\overline{\psi }(D_s,I_1,\mathrm{},I_n)e^{it[\frac{1}{L}(2\pi LD_s+H_\alpha )+𝒫_{k+1}]}(LD_s)^z|_{t=L}.$$
We can now give:
Proof of Theorem C: By (4.21) and the fact that $`e^{2\pi iLD_s}I`$ on $`L^2(S_L^1)`$ we get
$$a_{k\gamma }=\tau _{\gamma k}(\sqrt{\mathrm{\Delta }})=$$
$`(4.22)`$
$$Res_{z=0}Tr[\frac{1}{L}2(\pi LD_s+H_\alpha )+𝒫_{k+1}]^ke^{iH_\alpha }e^{iL𝒫_{k+1}}\overline{\psi }(\frac{I}{ϵD_s})(LD_s)^z$$
with
$$\overline{\psi }(\frac{I}{ϵD_s}):=\mathrm{\Pi }_{j=1}^n\overline{\psi }(\frac{I_j}{ϵD_s}).$$
From the well-known spectra of $`D_s,I_1,\mathrm{},I_n`$ we can rewrite (4.22) as
$`(4.23)`$
$$Res_{z=0}\underset{m=1}{\overset{\mathrm{}}{}}m^z\{\underset{(q_1,\mathrm{},q_n)𝐍^n}{}[\frac{1}{L}(2\pi m+\underset{j=1}{\overset{n}{}}(q_j+\frac{1}{2})\alpha _j)+𝒫_{k+1}(m,q_1+\frac{1}{2},\mathrm{},q_n+\frac{1}{2})]^k$$
$$e^{iL𝒫_{k+1}(m,q_1+\frac{1}{2},\mathrm{},q_n+\frac{1}{2},L)}\overline{\psi }(\frac{q}{ϵm})e^{i_{j=1}^n(q_j+\frac{1}{2})\alpha _j}\}.$$
Regarding
$`(4.24)`$
$$\underset{(q_1,\mathrm{},q_n)𝐍^n}{}\overline{\psi }(\frac{q}{ϵm})e^{i_{j=1}^n(q_j+\frac{1}{2})\alpha _j}$$
as a smooth function of $`\alpha ^n`$, we can further rewrite (4.23) as
$$Res_{z=0}\underset{m=1}{\overset{\mathrm{}}{}}m^z\{[\frac{1}{L}(2\pi m+\underset{j=1}{\overset{n}{}}\alpha _jD_{\alpha _j})+𝒫_{k+1}(m,D_{\alpha _1},\mathrm{},D_{\alpha _n},L)]^k$$
$`(4.25)`$
$$e^{iL𝒫_{k+1}(m,D_{\alpha _1},\mathrm{},D_{\alpha _n},L)}\underset{(q_1,\mathrm{},q_n)𝐍^n}{}\overline{\psi }(\frac{q}{ϵm})e^{i_{j=1}^n(q_j+\frac{1}{2})\alpha _j}\}.$$
Since $`\overline{\psi }(\frac{q}{ϵm})`$ is for each $`m`$ a finitely supported function of $`q`$, we can also rewrite (4.24) as
$$\underset{\delta 0}{lim}\mathrm{\Pi }_{j=1}^n\overline{\psi }(\frac{D_{\alpha _j}}{ϵm})\underset{q_j=0}{\overset{\mathrm{}}{}}e^{i_{j=1}^n\alpha _j(q_j+\frac{1}{2})\delta }$$
$`(4.26)`$
$$=\underset{\delta 0}{lim}\mathrm{\Pi }_{j=1}^n\overline{\psi }(\frac{D_{\alpha _j}}{ϵm})\mathrm{\Pi }_{j=1}^n\frac{e^{i\frac{1}{2}\alpha _j}}{1e^{i(\alpha _j+i\delta })}.$$
Recalling the definition of $`T(\alpha )`$ (0.12a), combining (4.23)-(4.26), and taking the limit as $`\delta 0`$, we get:
$$a_{\gamma k}=Res_{z=0}\underset{m=1}{\overset{\mathrm{}}{}}m^z\{[\frac{1}{L}(2\pi m+\underset{j=1}{\overset{n}{}}\alpha _jD_{\alpha _j})+𝒫_{k+1}(m,D_{\alpha _1},\mathrm{},D_{\alpha _n},L)]^k$$
$`(4.27)`$
$$e^{iL𝒫_{k+1}(m,D_{\alpha _1},\mathrm{},D_{\alpha _n},L)}\mathrm{\Pi }_{j=1}^n\overline{\psi }(\frac{D_{\alpha _j}}{ϵm})T(\alpha )\}.$$
Here, we use the fact that $`T(\alpha )`$ is a a temperate distribution on $`_\alpha ^n`$ with singular support on $`_{j=1}^n\times \mathrm{}\times 2\pi \mathrm{}\times ^n`$ (the factor of $``$ occurring in the jth position) and that the limit $`\delta 0`$ in (4.26) can be taken in $`𝒮^{}`$. Since the cut-off smooths out the singularity, the right side of (4.27) is a smooth function of $`\alpha `$ and can be evaluated at the special values of $`\alpha `$ determined by $`\gamma .`$
Since $`𝒫_{k+1}(m,D_{\alpha _1},\mathrm{},D_{\alpha _n},L)`$ is a symbol of order $`1`$ in $`m`$ with coefficients given by polynomials in the operators $`D_{\alpha _j}`$, we can expand the kth power in (4.27) as an operator-valued polyhomogeneous function of $`m`$. At least formally, we can also expand the exponential $`e^{iL𝒫_{k+1}(m,D_{\alpha _1},\mathrm{},D_{\alpha _n},L)}`$ in a power series and then expand each term in the power series as a polynomial in $`m^1`$. Collecting powers of m, the right side of (4.27) can be put, at least formally, in the form
$`(4.28)`$
$$Res_{z=0}\underset{m=1}{\overset{\mathrm{}}{}}\underset{j=o}{\overset{\mathrm{}}{}}m^{z+kj}_{k,kj}(D_\alpha ,L)\overline{\psi }(\frac{D_\alpha }{mϵ})T(\alpha ),$$
with $`_{k,kj}(D_\alpha ,L)`$ the coefficient of $`m^{kj}`$ in (4.27). To justify this manipulation, we have to deal with the remainder term in the Taylor expansion of the exponential. We have
$$e^{ix}=e_N(ix)+r_N(ix)$$
$$e_N(ix)=1+ix+\mathrm{}+\frac{(ix)^N}{N!}$$
$$r_N(ix)=(ix)^{N+1}b_N(ix)$$
$$b_N(ix)=_o^1\mathrm{}_o^1t_N^Nt_{N1}^{N1}\mathrm{}t_o^oe^{t_Nt_{N1}\mathrm{}t_oix}𝑑t_o\mathrm{}𝑑t_N.$$
Hence we have
$$e^{iL𝒫_{k+1}}:=e_N(iL𝒫_{k+1})+(iL𝒫_{k+1})^{N+1}b_N(iL𝒫_{k+1})$$
with $`𝒫_{k+1}`$ short for $`𝒫_{k+1}(m,D_{\alpha _1},\mathrm{},D_{\alpha _n})`$. The $`e_N`$ term of course contributes a finite number of terms of the desired form (4.28). For the remainder, we expand $`(iL𝒫_{k+1})^{N+1}`$ as a polynomial in $`m^1`$ with coefficients given by operators $`Q_{Np}(D_{\alpha _1},\mathrm{},D_{\alpha _n})`$ and observe that each term has a factor of $`m^{N1}`$. For each such term, we remove the coefficient operator $`Q_{Np}`$ from the sum $`_m`$, leaving only the factor of $`b_N`$. We then rewrite the resulting sum as a double sum $`_{mq}`$ as in (4.23), replacing all operators in $`D_{\alpha _j}`$ by their eigenvalues. Since $`b_N(ix)`$ is a bounded function and since each term of the resulting sum has at least the factor $`m^{zN1+k}`$ (possibly multiplied by a negative power of $`m`$), we see that the remainder is a sum of terms of the form
$`(4.29)`$
$$Res_{z=0}Q_{Np}(D_{\alpha _1},\mathrm{},D_{\alpha _n})\underset{kq}{}m^{z+kN1l}b_N(i𝒫_{k+1}(m,q+\frac{1}{2}))\overline{\psi }(\frac{q}{ϵm})e^{iq+\frac{1}{2},\alpha }.$$
We then observe that the sum is bounded by $`_{m=1}^{\mathrm{}}m^{RezN1+k+n}`$, hence converges absolutely and uniformly for $`Rez>N+k+n`$. It follows that for $`N>(n+k)`$ the sum in (4.29) defines a holomorphic function of $`z`$ in a half-plane containing $`z=0`$ and since the operations of taking the residue in $`z`$ and derivatives in $`\alpha `$ commute, each term (4.29) is zero. This justifies (4.28) and shows that it is actually a finite sum in j, say $`j<M`$ (in fact M=(k+1)(n+k+1)).
The residue in (4.28) is therefore well-defined and independent of $`ϵ`$. Since $`\overline{\psi }(\frac{D_\alpha }{mϵ})T(\alpha )T(\alpha )`$ as $`ϵ\mathrm{}`$ we must have
$`(4.30)`$
$$a_{\gamma k}=Res_{z=0}\underset{m=1}{\overset{\mathrm{}}{}}\underset{j=o}{\overset{M}{}}m^{z+kj}_{k,kj}(D_\alpha ,L)T(\alpha )$$
$$=Res_{z=0}\underset{j=0}{\overset{M}{}}\zeta (z+jk)_{k,kj}(D_\alpha ,L)T(\alpha ).$$
Here, $`\zeta `$ is the Riemann zeta-function, which has only a simple pole at $`s=1`$ with reside equal to one. It follows that the only term contributing to (4.29) is that with $`j=k+1`$ and hence we have
$`(4.30),`$
$$a_{\gamma k}=_{k,1}(D_\alpha ,L)T(\alpha )$$
completing the proof of Theorem C.∎
## 5. Local formulae for the residues: Proof of Theorem A
The characterization of the wave invariants in Theorem A is reminiscent of that of the heat invariants in \[ABP\] or \[Gi\], but involves non-local metric invariants near $`\gamma .`$ We begin by determining the metric data which contribute to $`a_{\gamma k}`$.
As always, we assume that $`\gamma `$ denotes a primitive closed geodesic, and denote the $`\mathrm{}`$th iterate of $`\gamma `$ by $`\gamma ^{\mathrm{}}`$. As above, we use the exponential map along $`N_\gamma S_L^1\times ^n`$ to pull back the metric to a metric on a neighborhood of $`S_L^1\times \{0\}`$, or more simply $`S^1`$, in $`S_L^1\times ^n`$ with the same wave invariants as $`g`$ along $`\gamma `$. This reduces the theorem to the case $`M=S_L^1\times ^n`$, with $`S^1`$ a closed geodesic of the metric.
We let $`J_{S^1}^m`$ denote the manifold of $`m`$-jets along $`S^1`$ of Riemannian metrics on $`S^1\times ^n`$ with the property that $`S^1`$ is a closed geodesic. We also let $`J_{S^1ell}^m`$ denote the open submanifold of m-jets of metrics for which $`S^1`$ is non-degenerate elliptic. We also write $`J_{\gamma ell}^m`$ when we wish to identify $`\gamma `$ and $`S^1`$. The density $`I_{\gamma k}(s)ds`$ of the $`k`$th wave invariant is then a map
$$I_{\gamma k}:J_{\gamma ell}^{m_k}\mathrm{\Omega }^1(S^1)$$
where $`\mathrm{\Omega }^1(S^1)`$ is the space of smooth densities along $`S^1`$ and where the jet order $`m_k`$ will be shown below to be $`m_k=2k+4`$. Since $`I_{\gamma k}(s,g)ds`$ is independent of the choice of coordinates on $`S^1\times ^n`$ we may express it in terms of Fermi normal coordinates $`(s,y)`$ with respect to a fixed normal frame $`e(s)`$ for g. As usual, the metric coefficients $`g_{ij}(i,j=o,1,\mathrm{},n)`$ will be understood relative to these coordinates.
§5.1 The metric data in $`I_{k\gamma }`$
We now claim that $`I_{k\gamma }`$ is an invariant polynomial in the following data:
(i) the curvature tensor $`R`$ and its covariant derivatives $`^mR`$ with $`m2k+2`$, contracted with respect to the Fermi normal vector fields $`\frac{}{s}`$ and $`_{j=1}^nc_j\frac{}{y_j}`$ with $`c_j`$;
(ii) the components $`Rey_{ij}(s),Imy_{ij}(s)`$ of the normalized eigenvectors $`Y_j𝒥_\gamma ^{}`$ of $`P_\gamma `$ relative to $`e_1,\mathrm{},e_n`$, and their first derivatives;
(iii) at most $`2k+1`$ indefinite integrals over $`S_L^1`$ of (i)-(ii).
(iv) the length $`L`$ and inverse length $`L^1`$ of $`\gamma ^{\mathrm{}}`$;
(v) the Floquet invariants $`\beta _j=(1e^{i\alpha _j})^1`$;
Indeed, by Theorem C, $`I_{\gamma ,k}(s,g)ds`$ is a density depending only on the data contained in $`T(\alpha ),L`$ and in the coefficients of $`_{k,1}`$. The latter is identical to the data in the coefficients in $`\stackrel{~}{p}_1,\mathrm{},\stackrel{~}{p}_{k+1}`$, hence to that in $`p_1,\mathrm{},p_{k+1}`$. By Proposition (4.2), these depends only on $`j_\gamma ^{2k+4}g`$, and on the coefficients of the symplectic matrices $`r_\alpha (s),a_s`$ which arise in the intertwining of $`\mathrm{\Delta }_h`$ to its normal form. The coefficients of $`r_\alpha `$ depend only on L and the $`\alpha _j`$’s and those of $`a_s`$ depend only on the Jacobi field components $`y_{ij},`$ and their first derivatives $`\dot{y}_{ij}`$. The second and higher derivatives of $`y_{ij}`$ can of course be eliminated by means of the Jacobi equation. Regarding (ii), we recall (see \[Gr, Theorem 9.21\]) that the Taylor series expansions at $`y=0`$ of the $`g_{ij}(s,y)`$’s, of the co-metric coefficients $`g^{ij}(s,y)`$, and of the volume densities and their powers, involve only the curvature tensor $`R`$ and its covariant derivatives $`R,\mathrm{},^mR`$ contracted with respect to the normal Fermi vector fields. Since the coefficients $`_{2\frac{m}{2}}`$ of the semi-classical expansion of $`\mathrm{\Delta }_h`$ depend only on this data, and since the intertwinings only introduce coefficients in data (ii)-(iv), the $`p_k(I)`$’s can only involve the $`g_{ij}`$’s thru the curvature and its covariant derivatives.
For the fact that $`I_{\gamma k}(s)`$ is a polynomial in (i)-(iv), it suffices to reconsider the construction of the normal form. Obviously the coefficients of the $``$’s and their metaplectic conjugates the $`𝒟_{\frac{j}{2}}`$’s are polynomials in (i)-(iv). It then follows from (2.39c) that the $`f_j(I)^{}s`$, hence the $`p_j(I)`$’s, are also polynomials in this data. Indeed, we argue inductively from the construction of the normal form that the coefficients change in the step from $`𝒟_{\frac{j}{2}}^{\frac{m}{2}}`$ to $`𝒟_{\frac{j}{2}}^{\frac{m+1}{2}}`$ as a result of taking commutators with operators whose coefficients are polynomials in (i)-(iv). It follows that the coefficients in $`𝒟_{\frac{j}{2}}^{\frac{m+1}{2}}`$ are also polynomials in this data, with possibly two more $`D_s`$-derivatives due to the commutators with $`D_s`$. Since the Jacobi data $`y_{ij},\dot{y}_{ij}`$ enters in thru a linear change of variables, the degree of the polynomial in this data will be closely related to the degree in the $`(y,\eta )`$ variables, which is twice the degree in the $`I`$ variables. The degree of the polynomial in the Jacobi data is however not necessarily the same as that (namely, k+2) in the $`I`$-variables: from the proof of Lemma (2.22), we see that the algorithm for computing the polynomials $`f_j(I_1,\mathrm{},I_n)`$ involves taking operator commutators (or Poisson brackets of symbols); this lowers the order in the $`I`$-variables but not in the Jacobi data which are coefficients of the polynomials in the $`I`$’s. We will show below (see ‘Jacobi degrees’) that the order in the Jacobi data is no more than $`6k+6`$.
§5.2 Weights of $`I_{k\gamma }`$ and of the data
To determine the weights of the polynomials in the data (i)-(iv), we now extend (and in part recall) the table in §1.4 describing how the various data transform under the metric rescaling $`gg_ϵ:=ϵ^2g.`$ As above, $`(s,y)`$ always refer to Fermi normal coordinates relative to $`g`$, and for notational simplicity we put $`s=y_o.`$ The notations $`,R,\mathrm{\Delta }`$ refer respectively to the Riemannian connection, curvature and Laplacian. We distinguish $`\mathrm{\Delta }`$ from the local expression $`\frac{1}{\sqrt{g}}_{i,j=o}^n_{y_i}g^{ij}\sqrt{g}_{y_j}`$ for $`\mathrm{\Delta }`$ relative to the Fermi normal coordinate frame.
| $`g`$ | $`ϵ^2g`$ |
| --- | --- |
| $`g_{ij}=g(_{y_i},_{y_j})`$ | $`g_{ij}=ϵ^2g(ϵ^1_{y_i},ϵ^1_{y_j})`$ |
| $`D_{s,y}^\alpha g_{ij}`$ | $`ϵ^{|\alpha |}D_{s,y}^\alpha g_{ij}`$ |
| $`,R,\mathrm{\Delta }`$ | $`,R,ϵ^2\mathrm{\Delta }`$ |
| $`\frac{1}{\sqrt{g}}_{i,j=o}^n_{y_i}g^{ij}\sqrt{g}_{y_j}`$ | $`ϵ^2\frac{1}{\sqrt{g}}_{i,j=o}^n_{y_i}g^{ij}\sqrt{g}_{y_j}`$ |
| $`L,(s,y)`$ | $`ϵL,(ϵs,ϵy)`$ |
| $`y_{ij}:=g(Y_i,e_j)`$ | $`ϵ^{\frac{1}{2}}y_{ij}=ϵ^2g(ϵ^{\frac{1}{2}}Y_j,ϵ^1e_j)`$ |
| $`K_{ij}=g(R(_s,e_i)_s,e_j)`$ | $`ϵ^2K_{ij}`$ |
The trace of the wave group thus scales as
$`(\mathrm{5.2.1})`$
$$Tre^{it\sqrt{\mathrm{\Delta }}}Tre^{i\frac{t}{ϵ}\sqrt{\mathrm{\Delta }}}$$
from which it follows that
$`(\mathrm{5.2.2}a)`$
$$a_{\gamma k}=res\sqrt{\mathrm{\Delta }}^ke^{iL\sqrt{\mathrm{\Delta }}}ϵ^ka_{\gamma k}.$$
This can also be seen from the fact that $`a_{\gamma k}`$ is the coefficient of $`(tL+i0)^klog(tL+i0)`$ which is homogeneous of degree k modulo smooth functions of $`t`$. Since $`a_{\gamma k}=_\gamma I_{\gamma k}(s)𝑑s`$, and the integral scales like $`ds`$, we also have
$`(\mathrm{5.2.2}b)`$
$$I_{\gamma k}(s,ϵ^2g)=ϵ^{k+1}I_{\gamma k}(s).$$
As a check on the normal form, let us verify that $`wgt(a_{\gamma k})=k`$ directly from a weight analysis of the normal form. It is obvious that $`=\frac{1}{L}(LD_s+H_\alpha )`$ scales like $`ϵ^1`$ and so does each term in the expansion
$$\sqrt{\mathrm{\Delta }}+\frac{p_1(I_1,\mathrm{},I_n)}{L}+\mathrm{}.$$
Moreover the $`p_k(I_1,\mathrm{},I_n)`$’s and $`𝒫_k`$’s are also of weight -1, as determined in Lemma (2.22) and Theorem B. Since the transition from $`p_\nu `$ to the $`\stackrel{~}{p}_{\nu +r}`$’s only involves multiplication of $`p_\nu `$ with the $`r`$th power of $`\frac{H_\alpha }{LD_s}`$ it is obvious that $`\stackrel{~}{p}_\nu `$ also has weight -1. Expanding the exponent in the residue calculation, we get that
$`(\mathrm{5.2.3})`$
$$a_{\gamma k}=\underset{N=0}{\overset{k+1}{}}\frac{i^N}{N!}res[\frac{1}{L}(LD_s+H_\alpha )+𝒫_k]^k(L𝒫_k)^Ne^{iH_\alpha }$$
with $`[\frac{1}{L}(LD_s+H_\alpha )+𝒫_k]`$ of weight -1 and with $`(L𝒫_k)`$ of weight 0. Hence $`wgt(a_{\gamma k})=k.`$
§5.3 Wave invariants and QBNF coefficients
They are related as follows: The terms in (5.2.3) contributing nontrivially to $`a_{\gamma k}`$ have the form
$`(\mathrm{5.3.1}a)`$
$$L^{k+j_1+\mathrm{}j_{k+1}}res(LD_s)^{\mathrm{}(j_1+2j_2+\mathrm{}(k+1)j_{k+1})}H_\alpha ^{j_o}\stackrel{~}{p}_1^{j_1}\stackrel{~}{p}_2^{j_2}\mathrm{}\stackrel{~}{p}_{k+1}^{j_{k+1}}e^{iH_\alpha }$$
with
$`(\mathrm{5.3.1}b)`$
$$j_o+j_1+\mathrm{}j_{k+1}+\mathrm{}=N+k,j_1+2j_2+\mathrm{}(k+1)j_{k+1}=\mathrm{}+1.$$
By the argument in the proof of Theorem C ( §4), taking the residue removes the factors of $`LD_s`$ and replaces an expression $`res(LD)^1F(I)e^{iH_\alpha }`$ by the value of $`F(D_\alpha +\frac{1}{2})T(\alpha )`$ at a regular point. It follows that
$`(\mathrm{5.3.2})`$
$$a_{\gamma k}=C_{k;j,\mathrm{}}L^{k+j_1+\mathrm{}j_{k+1}}[H_\alpha ^{j_o}\stackrel{~}{p}_1^{j_1}\stackrel{~}{p}_2^{j_2}\mathrm{}\stackrel{~}{p}_{k+1}^{j_{k+1}}](D_\alpha +\frac{1}{2})T(\alpha )$$
where the sum is taken over the indices specificed in (5.3.1b).
§5.4 Jacobi degrees
We now show that the degree of $`I_{\gamma k}`$ as a polynomial in the Jacobi data $`y_{ij},\dot{y}_{ij}`$ is $`6k+6.`$ The proof a detailed review of the construction of the polynomials $`f_j(I_1,\mathrm{},I_n)`$ in Lemma (2.22) as well as the relation between the Jacobi degrees of the polynomials $`f_j`$’s and of the operators $`_{k,1}`$ in Theorem C. For the sake of brevity, and since it is quite routine, we will be a little sketchy in a few of the details. We will use the notation $`Jdeg`$ for the degree of a polynomial in the metric data above with respect to the Jacobi field components.
(5.4.1) Lemma $`Jdeg\stackrel{~}{p}_{j+1}=Jdegp_{j+1}=Jdegf_j=6j+6.`$
Proof From the formulae (4.19) relating the $`\stackrel{~}{p}_j(I_1,\mathrm{},I_n)`$’s and the $`p_j(I_1,\mathrm{},I_n)`$, and from the fact that $`\frac{H_\alpha }{D_s}`$ has Jacobi degree 0, we see that
$`(\mathrm{5.4.2}a)`$
$$Jdeg\stackrel{~}{p}_j=\text{max}\{Jdegp_1,\mathrm{},Jdegp_j\}$$
On the other hand, the relation between the $`p_j`$’s and $`f_j`$’s is given by comparing (3.3) and (3.12):
$`(\mathrm{5.4.2}b)`$
$$f_j=\frac{2}{L}p_{j+1}+\underset{i+r+2=j+1}{}p_{i+1}p_{r+1}$$
Let us assume $`Jdegf_j=6j+6`$: an easy induction using (5.4.2b) then shows that $`Jdegp_{j+1}=6j+6`$, and another using (5.4.2a) shows that $`Jdeg\stackrel{~}{p}_{j+1}=6j+6.`$
Hence we must prove (5.4.2b). By (2.39c) we have that
$`(\mathrm{5.4.3})`$
$$Jdegf_j=Jdeg𝒟_j^{j+\frac{1}{2}}.$$
To determine the latter, we need to recall that $`𝒟_j^{j+\frac{1}{2}}`$ is the $`h^j`$-term in the expression
$$Ad(e^{ih^{j+\frac{1}{2}}Q_{j+\frac{1}{2}}})Ad(e^{ih^jQ_j})\mathrm{}Ad(e^{ih^{\frac{1}{2}}Q_{\frac{1}{2}}})𝒟_h$$
with $`Ad(V)A:=V^1AV`$. Using that $`Ad(e^{iB})=e^{iad(B)},`$ and expanding everything in a formal $`h`$-series, we find that $`𝒟_j^{j+\frac{1}{2}}`$ is the $`h^j`$-term in the series
$`(\mathrm{5.4.4})`$
$$\underset{n,N_{\frac{1}{2}},\mathrm{},N_{2j1}}{}h^{\frac{n}{2}2}(h^{j+\frac{1}{2}})^{N_{j+\frac{1}{2}}}(h^j)^{N_j}\mathrm{}(h^{\frac{1}{2}})^{N_{\frac{1}{2}}}ad(Q_{j+\frac{1}{2}})^{N_{j+\frac{1}{2}}}\mathrm{}ad(Q_{\frac{1}{2}})^{N_{\frac{1}{2}}}𝒟_{2\frac{n}{2}}.$$
We claim
(5.4.5) Claim: Suppose $`JdegQ_{\frac{m}{2}}=3m`$ for $`m2j+1`$. Then: $`Jdeg𝒟_j^{j+\frac{1}{2}}=6j+6`$.
Proof: The $`h^j`$ term in (5.4.4) is the sum of terms with indices satisfying $`\frac{n}{2}2+N_{j+\frac{1}{2}}(j+\frac{1}{2})+\mathrm{}+N_{\frac{1}{2}}(\frac{1}{2})=j.`$ Multiplying by 6 and using the hypothesis we find that
$$Jdeg𝒟_j^{j+\frac{1}{2}}=n+6\underset{m=1}{\overset{2j+1}{}}\frac{m}{2}N_{\frac{m}{2}}=n+6[j+2\frac{n}{2}].$$
The claim now follows from the fact that $`n3`$ in any term in the sum.∎
(5.4.6) Claim: Suppose $`Jdeg𝒟_{\frac{j}{2}+\frac{1}{2}}^{\frac{j}{2}}=3j+3.`$ Then: $`JdegQ_{\frac{j+1}{2}}=3j+3`$.
Proof: Follows from (2.39c,d) which implies that
$`(\mathrm{5.4.7})`$
$$JdegQ_{\frac{j}{2}+\frac{1}{2}}=Jdeg𝒟_{\frac{j}{2}+\frac{1}{2}}^{\frac{j}{2}}.$$
We then show:
(5.4.8) Claim: $`JdegQ_{\frac{j}{2}}=3j`$.
Proof: We prove this by induction on $`j`$. For $`j=1`$ it holds by explicit calculation: from (2.27a), $`degQ_{\frac{1}{2}}=deg𝒟_{\frac{1}{2}}=deg_{\frac{1}{2}}=3.`$ It follows then by Claim (5.4.5) (with j=0) that $`Jdeg𝒟_o^{\frac{1}{2}}=6`$, and then by Claim (5.4.6) (with j=1) that $`JdegQ_1=6.`$ The rest of the induction proceeds similarly and the details are left to the reader. ∎
The proof of Lemma (5.4.1) is completed by combining (5.4.3) and Claims (5..4.5) and (5.4.6).∎
Finally we have
(5.4.9) Lemma: $`Jdega_{k\gamma }6k+6`$.
Proof: As in §5.3, the kth residual terms in $`[(D_s+H_\alpha )+𝒫_k]^ke^{iH_\alpha }e^{iL𝒫_k}`$ involve only the products
$$\stackrel{~}{p}_1^{j_1}\stackrel{~}{p}_2^{j_2}\mathrm{}\stackrel{~}{p}_{k+1}^{j_{k+1}}$$
of the $`\stackrel{~}{p}_j`$’s with
$`(\mathrm{5.4.10})`$
$$j_o+j_1+\mathrm{}j_{k+1}+\mathrm{}=N+k,j_1+2j_2+\mathrm{}(k+1)j_{k+1}=\mathrm{}+1,\mathrm{}k,Nk+1.$$
Using that $`Jdeg\stackrel{~}{p}_m6m`$, the Jacobi degree of any such term cannot exceed $`6(j_1+\mathrm{}+j_{k+1}).`$ We claim that the maximum possible value, subject to the constraints (5.4.10), is 6k+6. Indeed, since $`j_1+\mathrm{}+j_{k+1}=\mathrm{}+1(j_2+2j_3+\mathrm{}kj_{k+1})k+1(j_2+2j_3+\mathrm{}kj_{k+1})`$ the maximum value is achieved when $`j_1=k+1`$ and all other $`j_m`$’s are zero.∎
This completes the proof of Theorem A. ∎
## 6. Quantum Birkhoff normal form coefficients
In this section we give a brief summary of the algorithm for calculating the quantum Birkhoff normal form coefficients $`B_{\gamma ;k,j}`$. From these coefficients one could determine the wave invariants $`a_{\gamma ;k}`$ as described in §5.3, but since the $`B_{\gamma ;k,j}`$’s are simpler spectral invariants we choose to concentrate on them. In the next section, we apply the algorithm to the calculation of the coefficients $`B_{\gamma ;o,j}`$ in dimension 2.
Preliminaries
(6.1) Definition The quantum Birkhoff normal form (QBNF) coefficients are the coefficients of the monomials $`I^j:=I_1^{j_1}\mathrm{}I_n^{j_n}=|z_1|^{2j_1}\mathrm{}|z_n|^{2j_n}`$ in the complete Weyl symbol $`\stackrel{~}{p}_k(|z|_1^1,\mathrm{},|z|_n^2)`$ of the coefficient operator $`\stackrel{~}{p}_k(I_1,\mathrm{},I_n)`$ of the model normal form of Theorem B:
$$\stackrel{~}{p}_k(|z|_1^1,\mathrm{},|z|_n^2)=\underset{j^n:|j|k+1}{}B_{\gamma ;k,j}|z|^{2j}.$$
There is of course something arbitrary in the emphasis on the $`\stackrel{~}{p}_k(I_1,\mathrm{},I_n)`$’s here. The coefficients of the monomials of the complete symbol of the $`p_k(I_1,\mathrm{},I_n)`$’s are equally spectral invariants and we only prefer the $`\stackrel{~}{p}_k(I_1,\mathrm{},I_n)`$’s to maintain consistency with the terminology of \[G.1-2\]. Moreover, the coefficients of the operator monomials $`I^j`$ in either the $`p_k`$’s or $`\stackrel{~}{p}`$’s are also spectral invariants and in view of the relation between wave invariants and the QBNF (§4), it is the operator coefficients which are most simply related to the wave invariants.
The crucial point for the effective calculation of the QBNF coefficients is that they can be obtained from the coefficients of the complete symbol of the semi-classical normal form (SCNF)
$$W_h^{}_hW_h|_oh^2+\underset{j=0}{\overset{\mathrm{}}{}}h^kf_k(I_1,\mathrm{},I_n)$$
of Lemma (2.22) (after restriction to weight 0.) Indeed, after substituting $`h=^1,`$ the square of the QBNF is the SCNF. Hence, the key invariants are really the coefficients of the monomials in the complete symbols
$$f_k(|z_1|^2,\mathrm{},|z_n|^2)=\underset{j^n,|j|k+2}{}c_{\gamma ;kj}|z|^{2j}$$
of the coefficient operators $`f_k(I_1,\mathrm{},I_n)`$. In the remainder of the section, we drop the subscript $`\gamma `$ from the notation.
On the operator level, the relation between the $`\stackrel{~}{p}_k`$’s, $`p_k`$’s and $`f_k^{}s`$ is very simple: the operators $`p_k(I_1,\mathrm{},I_n)`$ of Theorem B are related to the operators $`f_k(I_1,\mathrm{},I_n)`$ of Lemma (2.22) by:
$`(6.2a)`$
$$f_k=\frac{2}{L}p_{k+1}+\underset{i+j=k,i,j1}{}p_ip_j.$$
while the operators $`\stackrel{~}{p}_j`$ are related to the operators $`p_k`$ by:
$$\stackrel{~}{p}_{\mathrm{}}=\underset{k1,(j_1,\mathrm{},j_k):k+|j|=\mathrm{}}{}C_{k,(j_1,\mathrm{},j_k)}H_\alpha ^{|j|}p_k$$
for certain universal (multinomial) constants $`C_{k,(j_1,\mathrm{},j_k)}`$. On the symbol level, the relation is a little more complicated since one has to compose the symbols in the Weyl calculus. In the following we summarise the steps involved in calculating the coefficients of the monomials $`|z|^{2j}`$ in the complete symbols of the $`f_k`$’s. To distinguish on operator from its complete symbol we use the notation $`\widehat{f}(I_1,\mathrm{},I_n)`$ for the Weyl operator with complete symbol $`f(|z_1|^2,\mathrm{},|z_n|^2).`$ We also recall that the notation $`|_o`$ refers to the weight 0 part of an operator relative to $`D_s`$: $`(A_2(s,x,D_x)D_s^2+A_1(s,x,D_x)D_s+A_o(s,x,D_x))|_o=A_o(s,y,D_y).`$ We will assume that the $``$’s and $`𝒟`$’s have been expressed in terms of the weightless Fermi normal coordinats $`(s,x)`$ of §1.5. The complex coordinates $`z`$ are given by $`z=x+i\xi .`$
Summary of the algorithm
| Step | Relevant data |
| --- | --- |
| Step 1 | Express $`\widehat{f}_{\frac{N1}{2}}`$ in terms of $`𝒟_{2\frac{n}{2}}`$’s |
| Formula for $`\widehat{f}_{\frac{N1}{2}}`$ | $`\widehat{f}_{\frac{N1}{2}}(I_1,\mathrm{},I_n):=\frac{1}{L}_o^L_{T^n}V_t^{}𝒟_{\frac{N1}{2}}^{\frac{N}{2}}|_oV_tdtds`$ (diagonal part; N odd) |
| Equivalently | $`\widehat{f}_{\frac{N1}{2}}(q_1+\frac{1}{2},\mathrm{},q_n+\frac{1}{2}):=\frac{1}{L}_o^L𝒟_{\frac{N1}{2}}^{\frac{N}{2}}\gamma _q,\gamma _q𝑑s`$ |
| Weyl symbol | $`f_{\frac{N1}{2}}(|z_1|^2,\mathrm{},|z_n|^2)=\frac{1}{L}_{|k|\frac{N+3}{2}}c_{\frac{N1}{2};j}|z|^{2j}`$ |
| Recursion for $`𝒟_{\frac{N1}{2}}^{\frac{N+1}{2}}`$ | $`𝒟_{\frac{N1}{2}}^{\frac{N+1}{2}}=`$ term of order $`h^{\frac{N1}{2}}`$ in |
| $``$ | $`_{n=3}^{N+2}h^{2+\frac{n}{2}}𝒟_{2\frac{n}{2}}^{\frac{N}{2}}h^{\frac{N1}{2}}\{[L^1D_s,\stackrel{~}{Q}_{\frac{N+1}{2}}]+𝒟_{\frac{N1}{2}}^{\frac{N}{2}}\}`$ |
| $``$ | $`+_{n=N+4}^{\mathrm{}}h^{2+\frac{n}{2}}𝒟_{2\frac{n}{2}}^{\frac{N}{2}}+_{m=2}^{\mathrm{}}h^{m\frac{N+1}{2}}\frac{i^m}{m!}(ad\stackrel{~}{Q}_{\frac{N+1}{2}})^m(L^1D_s)`$ |
| $``$ | $`+_{m=1}^{\mathrm{}}_{p=3}^{\mathrm{}}h^{m\frac{N+1}{2}+\frac{p}{2}2}\frac{i^m}{m!}(ad\stackrel{~}{Q}_{\frac{N+1}{2}})^m(𝒟_{2\frac{p}{2}}^{\frac{N}{2}})`$ |
| Operator Transport : | $`N+1`$ odd: $`\{[L^1D_s,\stackrel{~}{Q}_{\frac{N+1}{2}}]+𝒟_{\frac{N1}{2}}^{\frac{N}{2}}\}|_o=0`$ |
| equations for $`\stackrel{~}{Q}`$’s | $`N+1`$ even: $`\{[L^1D_s,\stackrel{~}{Q}_{\frac{N+1}{2}}]+𝒟_{\frac{N1}{2}}^{\frac{N}{2}}\}|_o=f_{\frac{N1}{2}}(I_1,\mathrm{},I_n).`$ |
| Symbol | $`N+1`$ odd: $`L^1_s\stackrel{~}{Q}_{\frac{N+1}{2}}(s,z,\overline{z})=i𝒟_{\frac{N1}{2}}^{\frac{N}{2}}|_o(s,z\overline{z})`$ |
| Transport equations | $`N+1`$ even: $`L^1_s\stackrel{~}{Q}_{\frac{N+1}{2}}(s,z,\overline{z})=i(𝒟_{\frac{N1}{2}}^{\frac{N}{2}}\}|_o(s,z,\overline{z})+if_{\frac{N1}{2}}(|z_1|^2,\mathrm{},|z_1|^2).`$ |
| Homological | $`N+1`$ odd: $`\stackrel{~}{Q}_{\frac{N+1}{2}}(0,e^{i\alpha }z,e^{i\alpha }\overline{z})\stackrel{~}{Q}_{\frac{N+1}{2}}(0,z,\overline{z})=iL_0^L𝒟_{\frac{N1}{2}}^{\frac{N}{2}}|_o(s,z\overline{z})ds`$ |
| equations | $`N+1`$ even: $`\stackrel{~}{Q}_{\frac{N+1}{2}}(0,e^{i\alpha }z,e^{i\alpha }\overline{z})\stackrel{~}{Q}_{\frac{N+1}{2}}(0,z,\overline{z})=iL_0^L(𝒟_{\frac{N1}{2}}^{\frac{N}{2}}\}|_o(s,z,\overline{z})ds`$ |
| $``$ | $`f_{\frac{N1}{2}}(|z_1|^2,\mathrm{},|z_1|^2).`$ |
| Solutions: | N+1 odd: $`\stackrel{~}{Q}_{\frac{N+1}{2}}(s,z,\overline{z})=\stackrel{~}{Q}_{\frac{N+1}{2}}(0,z,\overline{z})iL_o^s𝒟_{\frac{N1}{2}}^{\frac{N}{2}}|_o(t,z,\overline{z})ds`$ |
| $``$ | N+1 even: $`\stackrel{~}{Q}_{\frac{N+1}{2}}(s,z,\overline{z})=\stackrel{~}{Q}_{\frac{N+1}{2}}(0,z,\overline{z})`$ |
| $``$ | $`iL_o^s\{𝒟_{\frac{N1}{2}}^{\frac{N}{2}}|_o(t,z,\overline{z})f_{\frac{N1}{2}}(|z_1|^2,\mathrm{},|z_n|^2)\}|_ods`$ |
| Step 2 | Express $`𝒟_{2\frac{n}{2}}`$’s in terms of metric data |
| Conjugate to $`_{2\frac{n}{2}}`$’s | $`𝒟_{2\frac{n}{2}}=\mu (𝒜_{}^{}{}_{}{}^{})_{2\frac{n}{2}}\mu (𝒜_{}^{}{}_{}{}^{})^1`$ |
| $`_{2\frac{n}{2}}`$ = | $`(hL)^2g_{[h]}^{oo}+2i(hL)^1g_{[h]}^{oo}_s+i(hL)^1\mathrm{\Gamma }_{[h]}^o`$ |
| $`h^{\frac{n}{2}2}`$-term of | $`+h^1(_{ij=1}^ng_{[h]}^{ij}_{x_i}_{x_j})+h^{\frac{1}{2}}(_{i=1}^n\mathrm{\Gamma }_{[h]}^i_{x_i})+(\sigma )_{[h]}`$ |
| Step 3 | Solve for coefficients |
| Coeff’s of $`𝒟_{2\frac{k}{2}}^{}|_o`$ | $`𝒟_{2\frac{n}{2}}^{}|_o(s,z,\overline{z})==_{|m|+|n|2k}d_{\frac{k}{2}2;m,n}^{}(s)z^m\overline{z}^n`$ |
| Coeff’s of $`𝒟_{\frac{N1}{2}}^{\frac{N}{2}}|_o`$ | $`d_{\frac{N1}{2};m,n}^{\frac{N}{2}}(s)=F_{\frac{N1}{2};m,n}^{\frac{N}{2}}(d_{\frac{k}{2}2;m,n};k\frac{N1}{2})`$ |
| Coeff’s of $`\stackrel{~}{Q}_{\frac{N+1}{2}}`$’s | $`\stackrel{~}{Q}_{\frac{N+1}{2}}(s,z,\overline{z})=_{|m|+|n|N+3}q_{\frac{N+1}{2};mn}(s)z^m\overline{z}^n`$ |
| Off-diagonal Coefficients | $`q_{\frac{N+1}{2};m,n}(0)=i(1e^{i(mn)\alpha })^1iL_o^Ld_{\frac{N1}{2};m,n}(s)𝑑s`$ |
| Diagonal Coefficients | $`c_{\frac{N1}{2};k}=\overline{d}_{\frac{N1}{2};k,k}^{\frac{N}{2}}=\frac{1}{L}_o^Ld_{\frac{N1}{2};k,k}^{\frac{N}{2}}(s)𝑑s`$ |
Remark Since $`\sqrt{\mathrm{\Delta }}`$ commutes with complex conjugation, the odd order terms in its Weyl symbol are even functions under the involution $`(s,\sigma ,x,\xi )(s,\sigma ,x,\xi ).`$ Although we have restricted to the positive cone where $`\sigma >0`$, one gets similar normal forms for $`\sigma <0`$ as long as $`D_s`$ is interpreted as $`|D_s|`$. Since the Harmonic oscillators and their powers have even symbols, it must be the case that the QBNF coefficients of the even order terms vanish. For instance the coefficient $`B_{o;2}`$ of $`\frac{I}{D_s}`$ must vanish automatically.
## 7. Explicit formulae in dimension 2
To illustrate the algorithm, we carry out the details of the calculation of $`a_{\gamma o}`$, or more importantly the normal form coefficients $`B_{k;j}`$ for k=0, in dimension 2. The result may be summarized as follows:
QBNF coefficients for k=0, dim =2The complete symbol of $`f_o(I)`$ in complex coordinates $`z=y+i\eta `$ has the form $`B_{o;4}|z|^4+B_{o;0}`$ where $`B_{o;j}`$ are given for both $`j=0,4`$ by weight -2 Fermi-Jacobi polynomials of the form:
$$B_{o;j}=\frac{1}{L}_o^L[a|\dot{Y}|^4+b_1\tau |\dot{Y}Y|^2+b_2\tau Re(\overline{Y}\dot{Y})^2+c\tau ^2|Y|^4+d\tau _{\nu \nu }|Y|^4+e\delta _{jo}\tau ]𝑑s+$$
$$+\frac{1}{L}\underset{0m,n3;m+n=3}{}C_{1;mn}\frac{sin((nm)\alpha )}{|(1e^{i(mn)\alpha })|^2}|_o^L\tau _\nu (s)\overline{Y}^mY^n](s)ds|^2$$
$$+\frac{1}{L}\underset{0m,n3;m+n=3}{}C_{2;mn}Im\{_o^L\tau _\nu (s)\overline{Y}^mY^n(s)[_o^s\tau _\nu (t)\overline{Y}^nY^m](t)dt]ds\}.$$
This corroborates the previous remark that the term $`B_{ok}`$ vanishes. As will be seen, the corresponding density is a total $`_s`$-derivative. Also, the coefficient $`\delta _{jo}`$ of $`\tau `$ is the Kronecker symbol, i.e. equals one if $`j=0`$ and vanishes if $`j=4`$.
The expressions for the normal form coefficients in higher dimensions are very similar, and it is only for the sake of simplicity that we have stated the result in dimension 2.
The result for the wave coefficient $`a_{\gamma o}`$ is then very similar, modulo the Floquet factors. Indeed, by (4.30) or by §5.3 the wave invariant is related to the normal form coefficients by
$$a_{\gamma o}=_{o,1}(D_\alpha ,L)T(\alpha )$$
with
$$_{o,1}(D_\alpha ,L)=L\stackrel{~}{p}_1(D_{\alpha _1}+\frac{1}{2},\mathrm{},D_{\alpha _n}+\frac{1}{2},L)=Lp_1(D_{\alpha _1}+\frac{1}{2},\mathrm{},D_{\alpha _n}+\frac{1}{2}).$$
Hence
$$_\gamma I_{\gamma o2}𝑑s=Lp_1(D_{\alpha _1}+\frac{1}{2},\mathrm{},D_{\alpha _n}+\frac{1}{2})T(\alpha )=Lp_1(D_{\alpha _1}+\frac{1}{2},\mathrm{},D_{\alpha _n}+\frac{1}{2})\mathrm{\Pi }_{j=1}^n(1e^{i\alpha _j})^1$$
In view of (2.33) and the fact that $`p_1=\frac{1}{2}Lf_o`$ we have
$`(7.1)`$
$$a_{\gamma o}=_\gamma I_{\gamma o}𝑑s=\frac{1}{2}L^2f_o(D_{\alpha _1}+\frac{1}{2},\mathrm{},D_{\alpha _n}+\frac{1}{2})\mathrm{\Pi }_{j=1}^n(1e^{i\alpha _j})^1.$$
From (7.1) and the formula for the complete symbol of $`f_o`$ one gets the explicit formula for $`a_{\gamma o}`$ as stated in the Introduction.
To prove that the $`B_{o;j}`$ coefficients have the form claimed above, we begin with the expression for $`f_o`$ from (2.33) or from the table in §6:
$`(7.2)`$
$$f_o(I_1,\mathrm{},I_n)=\frac{1}{L}_o^L_{T^n}V_t^{}𝒟_o^{\frac{1}{2}}V_t|_odt.$$
We also have, from (§2, (2.36) and below), that
$`(7.3)`$
$$𝒟_o^{\frac{1}{2}}=𝒟_o+\frac{1}{2}i[𝒟_{\frac{1}{2}},\stackrel{~}{Q}_{\frac{1}{2}}].$$
We wish to evalute the coefficients of the complete symbol of $`f_o`$ in terms of integrals over $`\gamma `$ of Fermi-Jacobi data. At first, we will allow the dimension to be arbitrary; when it is time to substitute in metric expressions we will restrict to dimension 2.
Let us consider first the second term on the right side, which simplifies the expression
$$[𝒟_{\frac{1}{2}},\stackrel{~}{Q}_{\frac{1}{2}}]+\frac{1}{2}i^2[[L^1D_s,\stackrel{~}{Q}_{\frac{1}{2}}],\stackrel{~}{Q}_{\frac{1}{2}}]$$
by using $`[L^1D_s,\stackrel{~}{Q}_{\frac{1}{2}}]=i𝒟_{\frac{1}{2}}.`$ We recall from §2 or from the table of §6 that the complete symbol of $`\stackrel{~}{Q}_{\frac{1}{2}}`$ is given (in complex vector notation) by
$$\stackrel{~}{Q}_{\frac{1}{2}}(s,z,\overline{z})=\stackrel{~}{Q}_{\frac{1}{2}}(0,z,\overline{z})+L_o^s𝒟_{\frac{1}{2}}|_o(t,z,\overline{z})dt,\stackrel{~}{Q}_{\frac{1}{2}}(0,e^{i\alpha }z,e^{i\alpha }\overline{z})\stackrel{~}{Q}_{\frac{1}{2}}(0,z,\overline{z})=L_o^L𝒟_{\frac{1}{2}}|_o(t,z,\overline{z})dt.$$
We note that $`𝒟_{\frac{1}{2}}`$ is independent of $`D_s`$ so that $`𝒟_{\frac{1}{2}}|_o=𝒟_{\frac{1}{2}}.`$ It follows that
$$[𝒟_{\frac{1}{2}}(s),\stackrel{~}{Q}_{\frac{1}{2}}(s)]=[𝒟_{\frac{1}{2}}(s),\stackrel{~}{Q}_{\frac{1}{2}}(0)]+[𝒟_{\frac{1}{2}}(s),_o^s𝒟_{\frac{1}{2}}(t)𝑑t]$$
so that the second term of (7.3) contributes to $`f_o(I_1,\mathrm{},I_n)`$ the diagonal part of
$`(7.4diag)`$
$$\frac{i}{2L^2}\{[\stackrel{~}{Q}_{\frac{1}{2}}(0,e^{i\alpha }z,e^{i\alpha }\overline{z}),\stackrel{~}{Q}_{\frac{1}{2}}(0,z,\overline{z}]+\frac{1}{2}\{[_o^L𝒟_{\frac{1}{2}}(s)ds,_o^s𝒟_{\frac{1}{2}}(t)dt].$$
Here, the bracket $`[,]`$ denotes the commutator of complete symbols in the sense of operator (or complete symbol) composition.
To determine the diagonal part, we write (as usual)
$$\stackrel{~}{Q}_{\frac{1}{2}}(0,z,\overline{z})=\underset{mn:|m|+|n|=3}{}q_{\frac{1}{2};mn}(0)z^m\overline{z}^n,𝒟_{\frac{1}{2}}(s,z,\overline{z})=\underset{mn:|m|+|n|=3}{}d_{\frac{1}{2};mn}(s)z^m\overline{z}^n.$$
Let us denote the operator composition of two complete Weyl symbols $`a,b`$ by $`a\mathrm{\#}b`$, that is
$$Op^w(a)Op^w(b)=Op^w(a\mathrm{\#}b)$$
and recall (cf. \[Ho III, Theorem 18.5.4\]) that this composition has an asymptotic expansion
$$a\mathrm{\#}bab+iP_1(a,b)+\frac{i^2}{2!}P_2(a,b)+\mathrm{}$$
where $`P_k(a,b)`$ is the higher order Poisson bracket (or transvectant), a bidifferential operator given in complex notation by
$$P_k(a,b)=(_z_{\overline{w}}_{\overline{z}}_w)^ka(z,\overline{z})b(w,\overline{w})|_{z=w}.$$
It is well known \[loc.cit.\] that the commutator is given symbolically by the odd expansion
$$a\mathrm{\#}bb\mathrm{\#}a\frac{1}{i}P_1(a,b)+\frac{1}{i^33!}P_3(a,b)+\mathrm{}$$
while anticommutators involve only the even transvectants. Since the complete symbols in (7.4 diag) are homogeneous polynomials of degree 3, the commutators involve only $`P_1`$ and $`P_3`$. One easily computes that
$$P_1(z^m\overline{z}^n,z^\mu \overline{z}^\nu )=C_{1;mn\mu \nu }z^{m+\mu 1}\overline{z}^{n+\nu 1}$$
where $`C_{1;mn\mu \nu }=\frac{1}{2}\sigma ((m,n),(\mu ,\nu ))`$ with $`\sigma `$ the standard symplectic inner product, and that
$$P_3(z^m\overline{z}^n,z^\mu \overline{z}^\nu )=C_{(m,n),(\mu ,\nu )}z^{m+\mu 3}\overline{z}^{n+\nu 3}$$
for certain other coefficients $`C_{(m,n),(\mu ,\nu )}`$. A straightforward computation then shows that the diagonal part in (7.4 diag) is the sum of the following terms:
$`(\mathrm{7.5.1})`$
$$\frac{i}{2L^2}\{[\stackrel{~}{Q}_{\frac{1}{2}}(0,e^{i\alpha }z,e^{i\alpha }\overline{z}),\stackrel{~}{Q}_{\frac{1}{2}}(0,z,\overline{z}]$$
$$=L^2[C_{1;3030}q_{\frac{1}{2};30}(0)q_{\frac{1}{2};03}(0)e^{i3\alpha }+C_{1;2112}q_{\frac{1}{2};21}(0)q_{\frac{1}{2};12}(0)e^{i2\alpha }+C_{1;1221}q_{\frac{1}{2};12}(0)q_{\frac{1}{2};21}(0)e^{i\alpha }+C_{1;0330}q_{\frac{1}{2};03}(0)q_{\frac{1}{2};30}(0)]|z|^4+$$
$$+[C_{3;3030}q_{\frac{1}{2};30}(0)q_{\frac{1}{2};03}(0)e^{i3\alpha }+C_{3;2112}q_{\frac{1}{2};21}(0)q_{\frac{1}{2};12}(0)e^{i2\alpha }+C_{3;1221}q_{\frac{1}{2};12}(0)q_{\frac{1}{2};21}(0)e^{i\alpha }+C_{3;0330}q_{\frac{1}{2};03}(0)q_{\frac{1}{2};30}(0)]$$
plus
$`(\mathrm{7.5.2})`$
$$\frac{1}{2}_o^L_o^s[\stackrel{~}{D}_{\frac{1}{2}}(s,z,\overline{z}),\stackrel{~}{D}_{\frac{1}{2}}(t,z,\overline{z}]dsdt$$
$$=_o^L_o^s[C_{1;3030}d_{\frac{1}{2};30}(s)d_{\frac{1}{2};03}(t)+C_{1;2112}d_{\frac{1}{2};21}(s)d_{\frac{1}{2};12}(t)+C_{1;1221}d_{\frac{1}{2};12}(s)d_{\frac{1}{2};21}(t)+C_{1;0330}d_{\frac{1}{2};30}(s)d_{\frac{1}{2};03}(t)dsdt]|z|^4+$$
$$_o^L_o^s[C_{3;3030}d_{\frac{1}{2};30}(s)d_{\frac{1}{2};03}(t)+C_{3;2112}d_{\frac{1}{2};21}(s)d_{\frac{1}{2};12}(t)+C_{3;1221}d_{\frac{1}{2};12}(s)d_{\frac{1}{2};21}(t)+C_{3;0330}d_{03}(s)d_{\frac{1}{2};30}(t)]𝑑s𝑑t.$$
We observe that there is no term of order $`|z|^2`$. Using that
$$q_{\frac{1}{2};mn}=iL(e^{i\alpha (mn)}1)^1_o^Ld_{\frac{1}{2};mn}(s)𝑑s$$
we reduce our problem to the calculation of the $`d_{\frac{1}{2};mn}(t)`$’s and the $`d_{o;mn}(t)`$’s.
To evaluate the expressions $`𝒟_{\frac{1}{2}}`$ and $`𝒟_o`$. we conjugate back to the $``$’s:
$$𝒟_o=\mu (𝒜_L^{})_o\mu ((𝒜_L^{})1,𝒟_{\frac{1}{2}}=\mu (𝒜_L^{})_{\frac{1}{2}}\mu (𝒜_L^{})^1(7.6a)$$
with
$`(7.6b)`$
$$𝒜_L:=\left(\begin{array}{cc}L^{\frac{1}{2}}Im\dot{Y}\hfill & L^{\frac{1}{2}}Re\dot{Y}\hfill \\ L^{\frac{1}{2}}ImY\hfill & L^{\frac{1}{2}}ReY\hfill \end{array}\right).$$
Using (2.11) we then compute the $``$’s in (unscaled) Fermi normal coordinates as follows:
$`(\mathrm{7.7.1}/2)`$
$$_{\frac{1}{2}}=\underset{\beta :|\beta |=3}{}\frac{1}{3!}_y^\beta g^{oo}(s,0)y^\beta $$
$`(\mathrm{7.7.1})`$
$$_o=\underset{\beta :|\beta |=4}{}\frac{1}{4!}_y^\beta g^{oo}(s,0)y^\beta \underset{\beta :|\beta |=2}{}_x^\beta [g^{oo}D_s+J^{\frac{1}{2}}D_s(g^{oo}J^{\frac{1}{2}})]|_{y=0}y^\beta $$
$$+\underset{\beta :|\beta |=2}{}\underset{ij=1}{\overset{n}{}}_y^\beta g^{ij}(s,0)y^\beta _{y_i}_{y_j}+\underset{i,j=1}{\overset{n}{}}_{y_j}\mathrm{\Gamma }^i(s,0)y_j_{y_i}$$
$$+_s^2+C_n\tau (s,0)$$
where $`\tau `$ is the scalar curvature. The last term comes from $`\mathrm{\Delta }_{\frac{1}{2}}1`$.
We then change variables $`y=Lx`$ and conjugate the symbols. By metaplectic covariance of the Weyl calculus, the conjugations change the complete Weyl symbols of the $``$’s (in the $`x`$ variables) by the linear symplectic transformation $`𝒜_L`$, i.e by the substitutions
$`(7.8)`$
$$\begin{array}{c}xL^{\frac{1}{2}}[(ReY)x+(ImY)\xi ]=\frac{1}{2}L^{\frac{1}{2}}[\overline{Y}z+Y\overline{z}]\hfill \\ \xi L^{\frac{1}{2}}[Re\dot{Y}x+(Im\dot{Y})\xi ]=\frac{1}{2}L^{\frac{1}{2}}[\overline{\dot{Y}}z+\dot{Y}\overline{z}]\hfill \end{array}.$$
It is evident that the normal form coefficients are going to be rather lengthy. To give the flavor of the full calculations in the simplest setting, we now specialize to the case of surfaces. Later we will briefly extend the calculations to all dimensions.
Dimension 2
In dimension 2 we have (in scaled Fermi coordinates)
$$g^{oo}(s,y)=1+C_1\tau (s)y^2+C_2\tau _\nu (s)y^3+\mathrm{}g^{11}=1J(s,u)=\sqrt{g_{oo}}=1+C_1^{}\tau (s)y^2+\mathrm{}$$
, for some constants $`C_j,C_j^{}`$ which will change from line to line. Hence the 1/2-density Laplacian in Fermi normal coordinates equals
$$\mathrm{\Delta }=J^{1/2}_sg^{oo}J_sJ^{1/2}+J^{1/2}_yJ_yJ^{1/2}$$
$$g^{oo}_s^2+\mathrm{\Gamma }^o_s+_y^2+\mathrm{\Gamma }^1_y+\sigma _o$$
and the rescaled Laplacian equals
$$\mathrm{\Delta }_h=(Lh)^2g_{[h]}^{oo}+2i(Lh)^1g_{[h]}^{oo}_s+i(Lh)^1\mathrm{\Gamma }_{[h]}^og_{[h]}^{oo}_s^2+\mathrm{\Gamma }_{[h]}^o_s+_y^2+h^{\frac{1}{2}}\mathrm{\Gamma }_{[h]}^1_y+(\sigma )_{[h]}$$
Using the Taylor expansion of the metric coefficients one finds that the $``$’s have the form
$`(7.9)`$
$$_{\frac{1}{2}}=CL^2\tau _\nu (s)y^3,_o=C_1L^2y^4\tau _{\nu \nu }+C_2L^1y^2\tau _s+C_3L^1\tau _sy^2_s^2+C_4\tau y_y+C_5\tau .$$
All terms have weight -2.
We now switch to the weightless coordinates $`y=Lx`$ and get:
$`(7.9)`$
$$_{\frac{1}{2}}=CL\tau _\nu (s)x^3,_o=C_1L^2x^4\tau _{\nu \nu }+C_2L^1x^2\tau _s+C_3L\tau _sx^2_s^2+C_4\tau x_x+C_5\tau .$$
Making the linear symplectic substitutions above we first get
$$𝒟_{\frac{1}{2}}(s,z,\overline{z})=CL^{\frac{1}{2}}\tau _\nu (s)([\overline{Y}z+Y\overline{z}])^3$$
hence
$`(\mathrm{7.10.1})`$
$$d_{\frac{1}{2};mn}(s)=C_{mn;3}L^{\frac{1}{2}}\tau _\nu [\overline{Y}^mY^n](s)$$
$`(\mathrm{7.10.2})`$
$$\overline{d}_{\frac{1}{2};mn}(s)=C_{mn;3}L^{\frac{3}{2}}_o^L\tau _\nu [\overline{Y}^mY^n](s)𝑑s$$
$`(\mathrm{7.10.3})`$
$$q_{\frac{1}{2};mn}=iC_{mn;3}(1e^{i(mn)\alpha })^1L^{\frac{1}{2}}_o^L\tau _\nu [\overline{Y}^mY^n](s)𝑑s$$
with $`m+n=3`$ and for certain coefficients $`C_{mn;\beta }`$. It is evident that $`d_{mn}(s)`$ is a Fermi-Jacobi polynomial of weight -2 and of Jacobi degree 3, so that the terms (7.5.1-2) are of Jacobi degree 6 as stated in Theorem A. The diagonal part has terms of degree $`|z|^4`$ and $`|z|^o`$ with coefficients of the form (with $`m,n=0,\mathrm{}3;m+n=3`$):
$$\frac{1}{L}\{[_o^L\tau _\nu (s)\overline{Y}^mY^n](s)ds][_o^L\tau _\nu (t)\overline{Y}^nY^m(t)dt]\frac{e^{i(nm)\alpha }}{|1e^{i(nm)\alpha }|^2}$$
$`(\mathrm{7.11.1})`$
$$[_o^L\tau _\nu (s)\overline{Y}^nY^m(s)ds][_o^L\tau _\nu (t)\overline{Y}^mY^n](t)dt]\frac{e^{i(mn)}}{|1e^{i(mn)}|^2}\}$$
and
$`(\mathrm{7.11.2})`$
$$\frac{1}{L}\{_o^L\tau _\nu (s)\overline{Y}^mY^n(s)[_o^s\tau _\nu (t)\overline{Y}^nY^m](t)dt]ds_o^L\tau _\nu (s)\overline{Y}^nY^m(s)[_o^s\tau _\nu (t)\overline{Y}^mY^n(t)dt]ds\}.$$
To calculate the complete symbol of $`D_s`$-weight 0, $`𝒟_o|_o`$, of the second term $`𝒟_o^{\frac{1}{2}}`$ we make the same linear substitution and eliminate any $`D_s`$ appearing all the way to the right. We also invert the relation
$$\mu (𝒜_L^{})^1D_s\mu (𝒜_L^{})=D_s\frac{1}{2}(L^1_x^2+L\tau x^2)$$
to get
$$\mu (𝒜^{})D_s\mu (𝒜^{})^1=D_s\frac{1}{2}\mu (𝒜^{})(L^1_x^2+L\tau x^2)\mu (𝒜^{})^1$$
and transform the complete symbol of quadratic term by the symplectic substitution. The result is that $`𝒟_o|_o(s,z,\overline{z})`$ equals
$`(\mathrm{7.12.1})`$
$$C_1\tau _{\nu \nu }[\overline{Y}z+Y\overline{z}]^4$$
$`(\mathrm{7.12.2})`$
$$+C_2L^2\tau [\overline{Y}z+Y\overline{z}]^2\mathrm{\#}([\overline{\dot{Y}}z+\dot{Y}\overline{z}]^2+\tau L^2[\overline{Y}z+Y\overline{z}]^2))$$
$`(\mathrm{7.12.3})`$
$$+C_3\tau _s[\overline{Y}z+Y\overline{z}]^2$$
$`(\mathrm{7.12.4.1})`$
$$2L^2_s([\overline{\dot{Y}}z+\dot{Y}\overline{z}]^2L^2\tau [\overline{Y}z+Y\overline{z}]^2)$$
$`(\mathrm{7.12.4.2})`$
$$+L^2\{[\overline{\dot{Y}}z+\dot{Y}\overline{z}]^2L^2\tau [\overline{Y}z+Y\overline{z}]^2\}\mathrm{\#}\{[\overline{\dot{Y}}z+\dot{Y}\overline{z}]^2L^2\tau [\overline{Y}z+Y\overline{z}]^2\}$$
$`(\mathrm{7.12.5})`$
$$+C_4L^1\tau (\overline{Y}z+Y\overline{z})\mathrm{\#}(\overline{\dot{Y}}z+\dot{Y}\overline{z})+C_5\tau .$$
Our concern is with the diagonal part of the complete symbol, that is, with the terms involving $`|z|^4,|z|^2,|z|^o`$, and more precisely with their integrals over $`\gamma `$. To begin with, we observe that the diagonal part of term (7.12.1) is purely of degree $`|z|^4`$ and its average over $`\gamma `$ equals
$`(\mathrm{7.13.1})`$
$$(Const.)|z|^4\frac{1}{L}_o^L\tau _{\nu \nu }|Y|^4ds.$$
The diagonal part of term (17.12.2) contributes only the $`P_o`$ and $`P_2`$ terms, of degrees $`|z|^4`$ and $`|z|^o`$ respectively, whose averages over $`\gamma `$ have the form
$`(\mathrm{7.13.2})`$
$$(|z|^4/or/|z|^o)\frac{1}{L}_o^L\tau [a\tau |Y|^4+bRe(\overline{Y}\dot{Y})^2+c|Y\dot{Y}|^2]𝑑s$$
with constants $`a,b,c,d`$ which can differ between the two degrees. The missing $`P_1`$-term vanishes: it is a multiple of the Poisson bracket
$$P_1([\overline{Y}z+Y\overline{z}]^2,\tau [\overline{Y}z+Y\overline{z}]^2))$$
which simplifies to a term of the form
$$\tau [\overline{Y}^2\dot{Y}^2Y^2\dot{\overline{Y}}^2]=\tau (\overline{Y}\dot{Y}Y\overline{\dot{Y}})(\overline{Y}\dot{Y}+Y\overline{\dot{Y}})=C\tau \frac{d}{ds}|Y|^2$$
by the symplectic normalization of the Jacobi eigenfield. However the integral over $`\gamma `$ of this term vanishes, that is
$`(\mathrm{7.13.3})`$
$$_o^L\tau _s|Y|^2=0.$$
This can be seen from the Jacobi equation, which implies:
$$[\overline{Y}(Y^{})^{\prime \prime }+\tau _s|Y|^2+\tau Y^{}\overline{Y}]=0;$$
integrating over $`\gamma `$ and integrating the first term by parts twice kills the outer terms and hence the inner one. In a similar way, the diagonal part of (7.12.3) is of pure degree $`|z|^2`$ with coefficient $`\tau _s|Y|^2`$, so it makes no contribution either. Nor does the term (7.12.4.1), which is manifestly a total derivative and hence automatically has zero integral. The term (7.12.4.2) is a composition square and hence contributes only a product $`P_o`$-term of degree $`|z|^4`$ and a $`P_2`$-term of degree $`|z|^o`$, namely (for j=0,2) the diagonal part of
$$P_j(z^2\dot{\overline{Y}}^2+2|z|^2|\dot{Y}|^2+\overline{z}^2\dot{Y}^2+\tau (z^2\overline{Y}^2+2|z|^2|Y|^2+\overline{z}^2Y^2,z^2\dot{\overline{Y}}^2+2|z|^2|\dot{Y}|^2+\overline{z}^2\dot{Y}^2+\tau (z^2\overline{Y}^2+2|z|^2|Y|^2+\overline{z}^2Y^2)$$
whose average over $`\gamma `$ has the form
$`(\mathrm{7.13.4.2})`$
$$(|z|^4/or/|z|^o)\frac{1}{L}_o^L[a|\dot{Y}|^4+b\tau Re(\dot{\overline{Y}}^2Y^2)+c\tau |\dot{Y}Y|^2+d\tau ^2|Y|^4]𝑑s$$
where again the coefficients may vary between the two degrees. Finally, we the first term of (7.12.5) obviously has no diagonal part while obviously the second term contributes
$`(\mathrm{7.13.5})`$
$$C\frac{1}{L}_o^L\tau 𝑑s.$$
This completes the analysis of the QBNF coefficients $`B_{o4},B_{o2},B_{o0}`$.
## 8. Appendix : The classical Birkhoff normal form
The method of §2 for putting $`\sqrt{\mathrm{\Delta }}`$ into quantum Birkhoff normal form began, essentially, by putting the linearization of $`\sqrt{\mathrm{\Delta }}`$ at $`\gamma `$ into normal form by a linear symplectic transformation, and then proceeded by induction on the jet filtration at $`\gamma `$. The purpose of this appendix is to describe, rather briefly and sketchily, how to put the principal symbol of $`\sqrt{\mathrm{\Delta }}`$ into Birkhoff normal form by an analogous method. (We have not found this particular algorithm in the literature, but it is quite likely that it, or a much better algorithm, is well-known). We hope that the classical algorithm will help clarify the procedure in the quantum case.
In the usual Fermi symplectic normal coordinates $`(s,\sigma ,y,\eta )`$, we may write the principal symbol of $`\sqrt{\mathrm{\Delta }}`$ in the form
$$\stackrel{~}{H}(s,\sigma ,y,\eta ):=(g^{oo}(s,y)\sigma ^2+\underset{ij=1}{\overset{n}{}}g^{ij}(s,y)\eta _i\eta _j)^{\frac{1}{2}}$$
$$=\sigma (g^{oo}(s,y)+\underset{ij=1}{\overset{n}{}}\frac{\eta _i\eta _j}{\sigma ^2})^{\frac{1}{2}}.$$
Taking the Taylor expansion at $`y=\eta =0`$ we get
$$\stackrel{~}{H}(s,\sigma ,y,\eta )=\sigma (1+\frac{1}{2}[\underset{ij=1}{\overset{n}{}}K_{ij}(s)y_iy_j+\underset{i=1}{\overset{n}{}}\frac{\eta _i^2}{\sigma ^2}]+\mathrm{})$$
from which we extract the linearized symbol
$`(A.1)`$
$$\stackrel{~}{h}(s,\sigma ,y,\eta )=\sigma +\frac{1}{2}\underset{ij=1}{\overset{n}{}}K_{ij}(s)y_iy_j\sigma +\underset{i=1}{\overset{n}{}}\frac{\eta _i^2}{\sigma }.$$
We make the symplectic change of variables
$`(A.2)`$
$$\varphi :(s,\sigma ,y,\eta )(s^{},\sigma ^{},y^{},\eta ^{}):=(s+\frac{1}{2}\frac{y\dot{\eta }}{\sigma },\sqrt{\sigma }y,\frac{\eta }{\sqrt{\sigma }})$$
which transforms $`\stackrel{~}{h}`$ into
$`(A.3)`$
$$h(s,\sigma ,y,\eta )=\sigma +\frac{1}{2}\underset{ij=1}{\overset{n}{}}K_{ij}(s)y_iy_j+\underset{i=1}{\overset{n}{}}\eta _i^2$$
and $`\stackrel{~}{H}`$ into
$$H(s,\sigma ,y,\eta )=h+h^{[3]}+\mathrm{}$$
with $`\mathrm{}`$ denoting terms of order 3 and higher in $`(y,\eta ).`$ (Such terms are of order 3/2 with respect to the order in the isotropic calculus, while $`h`$ is of order 1, which is the rationale for calling $`h`$ the linearized symbol).
The first step in putting $`H`$ into Birkhoff normal form is to put $`h`$ into Birkhoff normal form
$$\widehat{h}=\sigma +\frac{1}{2}\underset{i=1}{\overset{n}{}}\alpha _i(y_i^2+\eta _i^2)$$
by means of a symplectic map. Equivalently, we wish to convert the Hamilton equations
$$\begin{array}{c}\frac{d}{ds}s=1\hfill \\ \frac{d}{ds}\sigma =0\hfill \\ \frac{d}{ds}y=\eta \hfill \\ \frac{d}{ds}\eta =K(s)y\hfill \end{array}$$
into the linear equations
$$\begin{array}{c}\frac{d}{ds}s=1\hfill \\ \frac{d}{ds}\sigma =0\hfill \\ \frac{d}{ds}q=\alpha p\hfill \\ \frac{d}{ds}p=\alpha q.\hfill \end{array}$$
We first do this in just the $`(y,\eta ,q,p)`$ variables, with a symplectic map of the form $`(y,\eta )(q,p)=B(s)(y,\eta )`$. The condition on $`B`$ is that
$`(A.4)`$
$$\dot{B}B^1+B\stackrel{~}{K}B^1=\alpha \dot{J}$$
where $`\stackrel{~}{K}`$ is the block anti-diagonal matrix with blocks $`I`$ and $`K`$, where $`J`$ is the usual block anti-diagonal matrix with blocks $`\pm I`$ and where $`\alpha \dot{J}`$ is the block anti-diagonal matrix with coefficients $`\pm \alpha _j`$ replacing $`\pm 1`$ in $`J`$. The equation (A.4) is equivalent to
$`(A.5)`$
$$\frac{d}{ds}B^1+\stackrel{~}{K}(s)B^1\alpha \dot{B}^1J$$
which has the solution
$`(A.6)`$
$$B(s)=r_\alpha (s)a_s^1.$$
To make the map symplectic in all the $`(s,\sigma ,y,\eta )`$ variables, we observe that the map
$$\psi _1:(s,\sigma ,y,\eta )(s,\sigma \frac{1}{2}\underset{i=1}{\overset{n}{}}\alpha _i(y_i^2+\eta _i^2),r_\alpha (s)(y,\eta ))$$
is symplectic with respect to $`dsd\sigma +dyd\eta `$ and satisfies
$$\psi _1^{}(\sigma +\frac{1}{2}\underset{i=1}{\overset{n}{}}\alpha _i(y_i^2+\eta _i^2))=\sigma .$$
Also, the map
$$\psi _2:(s,\sigma ,y,\eta )(s,\sigma +f(s,y,\eta ),a_s^1(y,\eta ))$$
with $`f(s,y,\eta ):=\frac{1}{2}_{ij=1}^nK_{ij}(s)y_iy_j+_{i=1}^n\eta _i^2`$ is symplectic with respect to $`dsd\sigma +dyd\eta `$ and satisfies
$$\psi _2^{}(\sigma )=\sigma +f.$$
It follows that $`\chi _1:=\psi _2\psi _1`$ pulls back $`h`$ to its Birkhoff normal form.
Let us now write
$$H_1:=\chi _1^{}(H)=\widehat{h}+h_1^{[3]}+h_1^{[4]}+\mathrm{}$$
with $`h_1^{[k]}(s,\sigma ,y,\eta )`$ of order 1 and vanishing to order k at $`(y,\eta )=(0,0).`$ Following the algorithm for putting a Hamiltonian with non-degenerate minimum at $`0`$ into Birkhoff normal form \[AM, 5.6.8, p.500\] we seek a symplectic map of the form
$$\chi _2=expad\sigma ^{\frac{1}{2}}F_3$$
with $`F_3=F_3(s,y,\eta )`$ vanishing to order 3 and with
$`(A.7)`$
$$\{\widehat{h},[\chi _2^{}H_1]^3\}=0.$$
Here, $`\{,\}`$ denotes the Poisson bracket with respect to $`dsd\sigma +dyd\eta `$, $`[F]^k`$ denotes the terms of vanishing order $`k`$ and $`adFG:=\{F,G\}.`$ The equation (A.7) is equivalent to
$`(A.8)`$
$$h_1^{[3]}+\{\sigma ^{\frac{1}{2}}F_3,\widehat{h}\}ker(ad\widehat{h})𝒫^3$$
where $`𝒫^k`$ denotes the space of homogeneous polynomials of degree k in $`(y,\eta )`$ with smooth coefficients in $`s`$ and with an overall factor of $`\sigma ^{\frac{k}{2}1}`$. In terms of the complex coordinates $`z_j=y_j+i\eta _j`$ on $`^{2n}`$ we may write
$$F_3=\underset{|j|+|k|=3}{}c_{jk}(s)z^j\overline{z}^k$$
$$h_1^{[3]}=\underset{|j|+|k|=3}{}a_{jk}(s)z^j\overline{z}^k$$
with $`j=(j_1,\mathrm{},j_n),k=(k_1,\mathrm{},k_n)`$, and we may write the Lie derivative $``$ with respect to the Hamilton vector field of $`\widehat{h}`$ as
$$=\frac{}{s}+\underset{i=1}{\overset{n}{}}\alpha _i(z_i\frac{}{z_i}\overline{z}_i\frac{}{\overline{z}_i})$$
The monomials in $`kerad𝒫^k`$ are of the form $`z^j\overline{z}^k`$ with $`\alpha ,jk=0`$ which implies $`j=k`$ with our assumptions on $`\alpha .`$ No such terms occur for odd k, so equation (A.8) thus becomes
$`(A.9)`$
$$\underset{|j|+|k|=3}{}\dot{c}_{jk}(s)+\alpha (jk)c_{jk}(s)=a_{jk}(s).$$
Expanding the periodic (or more generally almost periodic) functions $`c_{jk}`$ and $`a_{jk}`$ in Fourier series
$$c_{jk}(s)=e^{ims}\widehat{c}_{jk}(m)a_{jk}(s)=e^{ims}\widehat{a}_{jk}(m)$$
we can solve (A.9) with
$`(A.10)`$
$$\widehat{c}_{jk}(m)=\frac{\widehat{a}_{jk}(m)}{im+\alpha ,jk}.$$
Writing $`H_2=\chi _2^{}\chi _1^{}H`$, we arrive at the analogous problem for the fourth order terms. As in the quantum case,the even steps behave a little differently from the odd ones since now there can be terms $`H_2^{[4]o}`$ in $`H_2^{[4]}`$ with $`|j|=|k|`$ and hence which lie in $`kerad𝒫^4.`$ Since these terms already commute with $``$ it suffices to solve the analogue of (A9) with only the coefficients $`a_{jk}(s)`$ coming from $`H_2^{[4]}H_2^{[4]o}`$.
This puts the terms up to fourth order in normal form, and the process continues inductively to define symplectic maps $`\chi _N\chi _{N1}\mathrm{}\chi _1`$ which pulls back H to a normal form up to degree N. Since $`\chi _N=ImodO_N`$ the infinite product defines a smooth symplectic map which pulls back $`H`$ to a normal form modulo $`O_{\mathrm{}}`$, that is, to its Birkhoff normal form.
## 9. Index of Notation
In the following, $`\tau _L`$ will denote the translation operator $`\tau _Lf(s,y)=f(s+L,y)`$ on functions on $`\times ^n`$. Operators $`A`$ then transform under $`\tau _L`$ by $`\tau _LA\tau _L^{}`$.
NI.1: Model objects
| Model object | $``$ |
| --- | --- |
| Functions | $`\tau _Lf=f`$ |
| Operators | $`\tau _LA\tau _L^{}=A`$ |
| Maximal abelian algebra | $`𝒜=<,I_1,\mathrm{},I_n>`$ |
| Distinguished element | $`:=\frac{1}{L}(LD_s+H_\alpha )`$ |
| Harmonic Oscillator | $`H_\alpha :=\frac{1}{2}_{k=1}^n\alpha _kI_k`$ |
| Gaussians/Hermites | $`\gamma _o(y)=\gamma _{iI}(q)=e^{\frac{1}{2}|y|^2}\gamma _q=C_qA_1^{q_1}\mathrm{}A_n^{q_n}`$ |
| $`𝒜`$-Eigenfunctions | $`\varphi _{kq}^o(s,y):=e_k(s)\gamma _q(y)`$ |
| $``$-Eigenvalues | $`\varphi _{kq}^o(s,y)=r_{kq}\varphi _{kq}^o(s,y),r_{kq}=\frac{1}{L}(2\pi k+_{j=1}^n(q_j+\frac{1}{2})\alpha _j)`$ |
| Scaled Laplacian | $`_h:=\mu (\stackrel{~}{a})\mathrm{\Delta }_h\mu (\stackrel{~}{a})^{}_{m=o}^{\mathrm{}}h^{(2+\frac{m}{2})}_{2\frac{m}{2}}`$ |
| Intertwiner to SC normal form | $`W_h:=\mathrm{\Pi }_{k=1}^{\mathrm{}}W_{h\frac{k}{2}}\mu (r_\alpha ),`$ $`\stackrel{~}{W}_{h\frac{k}{2}}:=exp(ih^{\frac{k}{2}}Q_{\frac{k}{2}})`$ |
| S.C. Normal form | $`W_h^{}_hW_h|_oh^2+_{j=0}^{\mathrm{}}h^kf_k(I_1,\mathrm{},I_n)`$ |
NI.2: Twisted model objects
| Twisted model object | $``$ |
| --- | --- |
| Functions | $`\tau _Lf=\mu (r_\alpha (L))f`$ |
| Operators | $`\tau _LA\tau _L^{}=\mu (r_\alpha (L))A\mu (r_\alpha (L))^{}`$ |
| Maximal abelian algebra | $`𝒜=<D_s,I_1,\mathrm{},I_n>`$ |
| Distinguished element | $`D_s`$ |
| Harmonic Oscillator | $`H_\alpha :=\frac{1}{2}_{k=1}^n\alpha _kI_k`$ |
| Gaussians/Hermites | $`\gamma _o(y)=\gamma _{iI}(q)=e^{\frac{1}{2}|y|^2}\gamma _q=C_qA_1^{q_1}\mathrm{}A_n^{q_n}`$ |
| $`𝒜`$-Eigenfunctions | $`e^{ir_{kq}s}\gamma _q(y)`$ |
| $`D_s`$-Eigenvalues | $`D_se^{ir_{kq}s}\gamma _q(y)=r_{kq}e^{ir_{kq}s}\gamma _q(y)`$ |
| Scaled Laplacian | $`𝒟_h:=\mu (a)\mathrm{\Delta }_h\mu (a)^{}_{m=o}^{\mathrm{}}h^{(2+\frac{m}{2})}𝒟_{2\frac{m}{2}}`$ |
| Intertwiner to SC normal form | $`\stackrel{~}{W}_h:=\mathrm{\Pi }_{k=1}^{\mathrm{}}\stackrel{~}{W}_{h\frac{k}{2}}\mu (r_\alpha ),`$ $`\stackrel{~}{W}_{h\frac{k}{2}}:=exp(ih^{\frac{k}{2}}\stackrel{~}{Q}_{\frac{k}{2}})`$ |
| S.C. Normal form | $`\stackrel{~}{W}_h^{}𝒟_h\stackrel{~}{W}_h|_oh^2+_{j=0}^{\mathrm{}}h^kf_k(I_1,\mathrm{},I_n)`$ |
NI.3: Adapted model objects
| Adapted model object | $``$ |
| --- | --- |
| Functions | $`\tau _Lf=\mu (T)f`$ |
| Operators | $`\tau _LA\tau _L^{}=\mu (T)A\mu (T)^{}`$ |
| Maximal abelian algebra | $`𝒜=<,\mathrm{\Lambda }_1\mathrm{\Lambda }_1^{},\mathrm{},\mathrm{\Lambda }_n\mathrm{\Lambda }_n^{}>`$ |
| Distinguished element | $`=D_s\frac{1}{2}(_{j=1}^nD_{u_j}^2+_{ij=1}^nK_{ij}(s)u_iu_j).`$ |
| Harmonic Oscillator | $`H_\alpha :=\frac{1}{2}_{k=1}^n\alpha _k\mathrm{\Lambda }_k\mathrm{\Lambda }_k^{}`$ |
| Gaussians/Hermites | $`\mu (a^1)\gamma _o(s,u):=U_o(s,u)=(detY(s))^{1/2}exp(\frac{i}{2}<\mathrm{\Gamma }(s)u,u>,`$ $`\mathrm{\Gamma }(s):=\frac{dY}{ds}Y^1;`$ |
| $``$ | $`\mu (a^1)\gamma _q:=U_q=C_q\mathrm{\Lambda }_1^{q_1}\mathrm{}\mathrm{\Lambda }_n^{q_n}U_o.`$ |
| $`𝒜`$-Eigenfunctions | $`\mu (\stackrel{~}{a}_s)(\varphi _{kq}^o)=\varphi _{kq}:=e^{ir_{kq}s}U_q(s,u)`$ |
| Scaled Laplacian | $`\mathrm{\Delta }_h_{m=0}^{\mathrm{}}h^{(2+m/2)}_{2m/2}`$ |
NI.4 Intertwining operators
1. $`\mu (r_\alpha ):=_{S^1}^{}\mu (r_\alpha (s))𝑑s`$
2. $`a_s:=\left(\begin{array}{cc}Im\dot{Y}(s)^{}\hfill & ImY(s)^{}\hfill \\ Re\dot{Y}(s)^{}\hfill & ReY(s)^{}\hfill \end{array}\right)=𝒜^{}.`$
3. $`𝒜_L(s):=\left(\begin{array}{cc}L^{\frac{1}{2}}Im\dot{Y}(s)\hfill & L^{\frac{1}{2}}Re\dot{Y}(s)\hfill \\ L^{\frac{1}{2}}ImY(s)\hfill & L^{\frac{1}{2}}ReY(s)\hfill \end{array}\right)`$
| Intertwining operator $`A`$ | $`\tau _LA\tau _L^{}`$ |
| --- | --- |
| $`\mu (a):`$ Adapted Model $``$ Twisted Model | $`\tau _L\mu (a)\tau _L^{}=\mu (r_\alpha )\mu (a)`$ |
| $`\mu (r_\alpha )`$ : Model $``$ Twisted Model | $`\tau _L\mu (r_\alpha )\tau _L^{}=\mu (r_\alpha (L))\mu (r_\alpha )`$ |
| With $`\stackrel{~}{a}_s:=a_sr_\alpha (s),\mu (\stackrel{~}{a})`$: Adpated Model $``$ Model | $`\tau _L\mu (\stackrel{~}{a})\tau _L^{}=\mu (\stackrel{~}{a})`$ | |
warning/0002/math0002094.html | ar5iv | text | # Parallel connections and bundles of arrangements revised July 17, 2000.
## 1 Introduction
Let $`V`$ be a vector space over a field $`𝕂`$. An arrangement $`𝒜`$ in $`V`$ is a finite collection of linear hyperplanes in $`V`$. The complement $`M=M(𝒜)`$ of $`𝒜`$ is $`V𝒜`$. A set of hyperplanes $``$ is dependent if the $`\mathrm{codim}()<||`$. These dependent sets determine a matroid $`G(𝒜)`$ with ground set $`𝒜`$, the underlying matroid of $`𝒜`$. Alternatively, $`G(𝒜)`$ is the linear matroid realized by the projective point configuration $`𝒜^{}`$ in $`(V^{})`$ determined by the defining linear forms for the hyperplanes of $`𝒜`$.
In case $`𝕂=`$ the complement $`M(𝒜)`$ is a connected manifold whose topology has been studied in great detail. In this case there is a strong connection between the topological structure of $`M(𝒜)`$ and the underlying matroid $`G(𝒜)`$. The paradigmatic result along these lines is that the cohomology of $`M(𝒜)`$ has a presentation depending only on $`G(𝒜)`$, with the consequence that the Poincaré series of the cohomology ring of $`M(𝒜)`$ essentially coincides with the characteristic polynomial of $`G(𝒜)`$ . It has become clear that techniques and constructions from matroid theory can have interesting and surprising implications for the topology of hyperplane complements. In this paper we interpret the matroidal notions of modular flat, principal truncation, and generalized parallel connection in this vein, in terms of bundles of complex hyperplane arrangements, their fibers, and pullbacks via inclusion maps.
Henceforth we restrict our study to complex arrangements. The intersection lattice $`L=L(𝒜)`$ of $`𝒜`$ is the set of subspaces $`X`$ of $`^{\mathrm{}}`$ which are intersections of hyperplanes of $`𝒜`$, $`X=`$ for $`𝒜`$, partially ordered by reverse inclusion. The smallest element of $`L`$ is $`O_L=^{\mathrm{}}`$, the empty intersection, and the largest element of $`L`$ is $`1_L=𝒜`$. For $`X,YL(𝒜),`$ the join $`XY`$ is $`XY`$ and the meet $`XY`$ is $`\{H𝒜|HX+Y\}`$. The rank function $`r`$ of $`L`$ is given by $`r(X)=\mathrm{codim}(X)`$, and the semimodular law holds:
$$r(XY)+r(XY)r(X)+r(Y)$$
for $`X,YL`$. Then $`L`$ is a geometric lattice, isomorphic to the lattice of flats of the matroid $`G(𝒜)`$, via the identification of $`XL`$ with the flat
$$𝒜_X=\{H𝒜|HX\}.$$
We will often refer to elements $`XL(𝒜)`$ as flats, tacitly identifying $`X`$ with $`𝒜_X`$. For instance, “point” and “line” refer to flats of rank one and two. The corank of a flat $`X`$ is $`r(1_L)r(X)`$, and “copoints” and “colines” are flats of corank one and two.
When equality holds in the formula above, $`(X,Y)`$ is called a modular pair. An element $`XL(𝒜)`$ is modular if $`(X,Y)`$ is a modular pair for every $`YL(𝒜)`$. This is equivalent to the condition that $`X+Y`$ be an element of $`L`$ for every $`YL`$. Let $`\pi `$ be the linear projection of $`^{\mathrm{}}`$ onto the quotient $`^{\mathrm{}}/X`$. Modularity of $`X`$ implies that fibers of $`\pi `$, the parallel translates of $`X`$, intersect each $`YL(𝒜)`$ in the same way, independent of position. This observation was already made by Terao in , who proved that $`\pi |_{M(𝒜)}`$ is a fiber bundle projection in case $`X`$ has corank one. But, in fact, it is easy to show that modularity of $`X`$ is equivalent to $`\pi `$ being a map of stratified spaces, under the natural stratifications of $`^{\mathrm{}}`$ and $`^{\mathrm{}}/X`$ determined by $`𝒜`$ and $`𝒜_X`$. Being a linear projection, it is trivial to show $`\pi `$ restricts to a submersion on each stratum. L. Paris showed how to extend $`\pi `$ to a proper map of stratified spaces. Then Thom’s Isotopy Lemma implies that $`\pi |_{M(𝒜)}`$ is a fiber bundle projection for $`X`$ a modular flat of arbitrary rank.
This fibration result interpolates between two well-known extreme cases. In case $`X`$ is a modular copoint, the result was proven in , as already mentioned. This case gives rise to the notion of supersolvable arrangement, and its connection with fiber-type arrangements , a much-studied class . In case $`X`$ is a point, i.e., a hyperplane of $`𝒜`$, then $`X`$ is automatically modular, and the fibration is just the restriction of the defining form $`\varphi :^{\mathrm{}}`$ of the hyperplane $`X`$. This gives rise to the well-known elementary “cone-decone” construction . The restriction $`\varphi |_{M(𝒜)}:M(𝒜)^{}`$ is in fact a trivial fibration, with fiber isomorphic to the complement in $`^\mathrm{}1`$ of an affine arrangement, the decone of $`𝒜`$.
The general modular fibration theorem was proved by L. Paris . At the same time, we were independently conducting the research reported on in this paper , and had arrived at the same conclusion, only to later discover an error in our treatment of the proper extension of $`\pi `$. We sketch the argument here, and refer the reader to for a complete proof, concentrating instead on other structural results and consequences of the theorem.
In , Terao establishes the result for modular copoints, and proves that for general modular $`X`$ the fibers of $`\pi `$ have the same combinatorial type. But he specifically remarks that a proof of local triviality in the general case is not at hand. See Remark 2.7. The proof of Corollary 3.2 is in a sense a parametrized version of the argument of , where stratification techniques were first used in the theory of arrangements, several years after Terao’s work.
The characteristic polynomial of a lattice was defined by G.-C. Rota. The characteristic polynomial of a matroid is the characteristic polynomial of its lattice of flats. The modular flat $`X`$ gives rise to a factorization of the characteristic polynomial of $`G(𝒜)`$ over the integers, with one factor given by the characteristic polynomial of $`G(𝒜_X)`$. This is Stanley’s modular factorization theorem . Brylawski identified the other factor as the characteristic polynomial of a related matroid, the complete principal truncation $`\overline{T}_X(G)`$ \[36, Section 7.4\] of $`G=G(𝒜)`$ along $`X`$, divided by $`(t1)`$. The complete principal truncation is obtained by successively adjoining generic points on the specified flat and contracting on the new points. Technically this is a matroid with multiple points; when we refer to $`\overline{T}_X(G)`$ we will always mean the associated simple matroid (with the same lattice of flats).
We show in Theorem 2.1 that the fiber of the bundle map $`\pi |_{M(𝒜)}`$ is the complement of the decone of an arrangement realizing the complete principal truncation of $`G(𝒜)`$ on the flat $`X`$. In addition, just as in the corank-one case , the monodromy of the bundle is shown to act trivially on the cohomology of the fiber (Theorem 3.5). Then the $`E_2`$ term in the Leray-Serre spectral sequence of $`\pi |_{M(𝒜)}`$ is isomorphic to the tensor product of the cohomology of the base $`M(𝒜_X)`$ with that of the fiber. Using the identity relating characteristic polynomials and Poincaré polynomials, we obtain a topological interpretation of the Stanley-Brylawski and Terao factorization results. In fact, the factorization of the characteristic polynomial implies that the spectral sequence degenerates at the $`E_2`$ term, just as in the corank-one case, although we have no topological proof of this fact (Remark 3.10).
In the corank-one situation, the monodromy of the bundle gives rise to a “braid monodromy” homomorphism from $`\pi _1(M)`$ to the (pure) braid group on $`n`$ strands, where $`n=|𝒜𝒜_X|`$. In the general case the analogue of this braid monodromy takes values in the fundamental group of the matroid stratum of the Grassmannian, or equivalently, the projective realization space, of the complete principal truncation $`\overline{T}_X(G)`$. See Remark 3.4.
The current research grew out of an attempt to clarify and generalize the construction of , which involved arrangements whose matroids are parallel connections. We began by studying the matroidal notion of generalized parallel connection. Loosely speaking, this is the free sum of two matroids along a common flat. This free sum is well-defined if and only if the flat is modular in one of the matroids. Thus we were led to the consideration of modular flats. The combinatorial study of Sections 2 and 4 formed the main part of an NSF Research Experiences for Undergraduates project in the summer of 1997. This work was reported on in , which provided the groundwork for this paper.
Given the modular fibration result, we show that generalized parallel connection, in a natural realization in terms of complex arrangements, corresponds to the pullback of fiber bundles (Theorem 4.2). The construction of , which yields diffeomorphisms of the complements of arrangements with non-isomorphic matroids, uses ordinary parallel connection, in which the identified flats are points. Then the diffeomorphisms of are a consequence of two elementary observations, that the cone-decone construction yields a trivial bundle, and that the pullback of a trivial bundle is trivial.
When the base and fiber of a modular fibration are both aspherical, it follows that the complement $`M(𝒜)`$ is also aspherical. In this case $`𝒜`$ is called a $`K(\pi ,1)`$ arrangement. The problem of identifying $`K(\pi ,1)`$ arrangements has been an important one in the study of complex arrangements. There are two well-known classes of $`K(\pi ,1)`$ arrangements, the supersolvable ones, which are abundant in all ranks, and the simplicial ones, which are rare in ranks greater than three. Other techniques for identifying $`K(\pi ,1)`$ arrangements are mostly restricted to arrangements of rank three. See for further exposition of the $`K(\pi ,1)`$ problem.
Corollary 3.2 provides a method for identifying $`K(\pi ,1)`$ arrangements in ranks greater than three. In the final section we exhibit new families of such arrangements, arising from the work of P. Edelman and V. Reiner and D. Bailey on threshold graphs and subarrangements of the Coxeter arrangement of type $`B_{\mathrm{}}`$. By the classification result of Bailey, these new examples are all “factored” \[23, Section 3.3\]. So our result provides more evidence for the conjecture that factored arrangements of arbitrary rank are $`K(\pi ,1)`$ . We also give an example of an arrangement of rank four which has two different modular colines. Then Corollary 3.2 implies that a certain arrangement of rank three, the cone of one of the fibers, is not $`K(\pi ,1)`$, an arrangement to which existing techniques do not apply.
The search for examples of high-rank $`K(\pi ,1)`$ arrangements was motivated by a suggestion of G. Ziegler several years ago concerning counterexamples to the “homotopy type conjecture,” that complex arrangements with the same underlying matroid should have homotopy-equivalent complements . This idea, laid out in , is to find $`K(\pi ,1)`$ arrangements of high rank, whose underlying matroids have different characteristic polynomials, but have isomorphic generic rank-three truncations. Then generic 3-dimensional sections of these arrangements will have isomorphic underlying matroids, but non-isomorphic fundamental groups. Unfortunately we are so far unable to construct such examples using the technique of this paper.
## 2 Projections and principal truncations
In this section we establish terminology and analyze the combinatorics associated with projections of hyperplane arrangements. Let $`𝒜`$ be a central arrangement of hyperplanes in $`^{\mathrm{}}`$. Let $`L=L(𝒜)`$ be the intersection lattice of $`𝒜`$, consisting of subspaces of $`^{\mathrm{}}`$ as described in the introduction.
Let $`XL`$. Let $`𝒜_X=\{H𝒜|HX\}`$, and let $`\pi :^{\mathrm{}}^{\mathrm{}}/X`$ be the natural projection. Note that $`\pi `$ maps each hyperplane $`H𝒜_X`$ to a hyperplane of $`^{\mathrm{}}/X`$. Henceforth we consider the arrangement $`𝒜_X`$ to be an arrangement in $`^{\mathrm{}}/X`$.
We shall have occasion to study arrangements formed by the intersections of the hyperplanes of $`𝒜`$ with a given affine subspace $`S`$. This induced arrangement in $`S`$ is called the restriction of $`𝒜`$ to $`S`$, denoted $`𝒜^S`$.
We start by describing the combinatorial structure of the affine arrangement $`𝒜_\pi `$ formed by restricting $`𝒜`$ to a generic fiber of $`\pi `$. This requires some discussion of the cone-decone construction of , and a description of the matroid construction of principal truncation along a flat.
There is a natural correspondence between arrangements of linear hyperplanes in $`^{\mathrm{}}`$ and arrangements of affine hyperplanes in $`^\mathrm{}1`$. The analytic operations need not concern us here; they are described in detail in and . One places a copy of a given affine $`(\mathrm{}1)`$-arrangement $`𝒜`$ into $`\{1\}\times ^\mathrm{}1^{\mathrm{}}`$. Then replace each of the affine subspaces in this copy of $``$ by its linear span in $`^{\mathrm{}}`$ (i.e., “cone over the origin”) and adjoin the “hyperplane at infinity” $`\{0\}\times ^\mathrm{}1`$, to obtain a central arrangement $`c𝒜`$, the cone of $`𝒜`$, in $`^{\mathrm{}}`$. The inverse operation, called “deconing,” takes a central arrangement $`𝒜`$ to its projective image, and dehomogenizes relative to a hyperplane $`H_{\mathrm{}}𝒜`$ to obtain an affine $`(\mathrm{}1)`$-arrangement $`d𝒜`$.
The intersections of hyperplanes of $`d𝒜`$ form a geometric semilattice , isomorphic to a subposet of $`L(𝒜)`$. Specifically,
$$L(d𝒜)\{YL(𝒜)|YH_{\mathrm{}}\}.$$
The fiber arrangement $`𝒜_\pi `$ is an affine arrangement of dimension $`\mathrm{}r(X)`$. We will show that the underlying matroid of the cone $`c𝒜_\pi `$ is the complete principal truncation of the matroid $`G(𝒜)`$ along the flat $`𝒜_X`$.
The principal truncation $`T_F(G)`$ of a matroid $`G`$ along a flat $`F`$ is constructed by adding a generic point $`p`$ on the flat $`F`$ and then contracting $`G`$ on $`p`$ \[36, Section 7.4\]. The result may be a matroid with multiple points. We tacitly simplify the resulting matroid, by removing any multiple points. This does not affect the intersection lattice, characteristic polynomial, or Orlik-Solomon algebra.
This operation can be iterated. The complete principal truncation $`\overline{T}_F(G)`$ is the result of $`r(F)1`$ successive principal truncations on $`F`$, so that $`F`$ reduces to a point. Equivalently, one can add $`r(F)1`$ generic points to $`F`$ and contract $`G`$ on the flat spanned by the new points. Contraction of a matroid on a point corresponds to projection of a projective point configuration from one of its points, or restriction of a hyperplane arrangement to one of its hyperplanes.
###### Theorem 2.1
Let $`XL`$ and let $`𝒜_\pi `$ be the affine arrangement obtained by restricting $`𝒜`$ to a generic fiber of $`\pi :^{\mathrm{}}^{\mathrm{}}/X`$. Then the matroid $`G(c𝒜_\pi )`$ is isomorphic to the complete principal truncation $`\overline{T}_X(G)`$ of $`G=G(𝒜)`$ along the flat $`𝒜_X`$.
###### Proof.
Dualizing the description of complete principal truncation to hyperplane arrangements, we see that $`\overline{T}_X(G)`$ is the matroid of the arrangement $`𝒜^P`$ obtained by choosing a generic subspace $`P`$ of codimension $`r(X)1`$ containing $`X`$, and restricting $`𝒜`$ to $`P`$. Then $`P`$ has dimension $`dim(X)+1`$, $`XP`$ is a hyperplane of $`𝒜^P`$, and an affine translate of $`XP`$ is a generic fiber of $`\pi `$. It follows that $`d(𝒜^P)𝒜_\pi `$, so $`𝒜^Pc𝒜_\pi `$. ∎
###### Definition 2.2
A pair $`(X,Y)`$ forms a modular pair in $`L`$ if
$$r(XY)+r(XY)=r(X)+r(Y).$$
An element $`XL`$ is modular if $`(X,Y)`$ is a modular pair for every $`YL`$.
The following lemma is the key to the proof of the modular fibration theorem, and is trivial to prove.
###### Lemma 2.3
Let $`X,YL`$. Then $`(X,Y)`$ is a modular pair if and only if $`X+YL`$. ∎
When $`X`$ is modular, the conclusion of Theorem 2.1 holds for every fiber of $`\pi `$ over points not in $`𝒜_X`$. To prove this we need to describe the rank function $`r_T`$ on the lattice of flats $`L(\overline{T}_X(G))`$. According to \[36, Proposition 7.4.9\], the set $`L(\overline{T}_X(G))`$ can be identified with $`\{YL|XY=0_{L(𝒜)}\text{or}YX\}`$. With this identification the rank function $`r_T`$ is given by
$$r_T(Y):=\{\begin{array}{cc}r(Y)\text{ if}XY=0_{L(𝒜)},\hfill & \\ r(Y)r(X)+1\text{if}YX.\hfill & \end{array}$$
###### Theorem 2.4
Suppose $`X`$ be a modular flat. Let $`\overline{v}=v+X(^{\mathrm{}}/X)𝒜_X`$ and let $`𝒜_{\overline{v}}`$ be the restriction of $`𝒜`$ to $`\pi ^1(\overline{v})`$. Then the intersection lattice $`L(c𝒜_{\overline{v}})`$ is isomorphic to $`L(\overline{T}_X(M)).`$
###### Proof.
As in the proof of Theorem 2.1, the arrangement $`c𝒜_{\overline{v}}`$ can be identified with the restriction $`𝒜^P`$ of $`𝒜`$ to the linear subspace $`P`$ of codimension $`r(X)1`$ spanned by $`X`$ and $`v`$. Then $`L(c𝒜_{\overline{v}})=\{PY|YL(𝒜)\}.`$ There are three cases.
* Case 1. Suppose $`YL(𝒜)`$ satisfies $`XY=0_{L(𝒜)}`$. By modularity of $`X,`$ $`X+Y=^{\mathrm{}}`$. Then there exists $`yY`$ such that $`\overline{v}=\overline{y}=y+X`$. Then $`PY=(v+X)Y=(y+X)Y=y+(XY)`$. Since $`\overline{v}𝒜_X`$, $`yX`$, so
$$\begin{array}{cc}\hfill \mathrm{codim}_P(PY)& =\mathrm{codim}_{y+X}(y+(XY))\hfill \\ & =\mathrm{codim}_X(XY)\hfill \\ & =\mathrm{codim}_{^{\mathrm{}}}(XY)\mathrm{codim}_{^{\mathrm{}}}(X)\hfill \\ & =r(XY)r(X)\hfill \\ & =r(Y),\hfill \end{array}$$
the last equality by modularity of $`X`$.
* Case 2. Suppose $`YX`$. Then $`YXP`$ so
$$\begin{array}{cc}\hfill \mathrm{codim}_P(PY)& =\mathrm{codim}_P(Y)\hfill \\ & =\mathrm{codim}_X(Y)+\mathrm{codim}_P(X)\hfill \\ & =r(Y)r(X)+1.\hfill \end{array}$$
* Case 3. Suppose $`0_{L(𝒜)}<XY<X`$. Then $`XYH`$ for some $`H𝒜`$. Note that $`vH`$, since $`\overline{v}𝒜_X`$. It follows that $`PH=(v+X)H=X`$ since $`XH`$ while $`vH`$. Since $`PYPH`$ we have $`PY=PY^{}`$ for $`Y^{}=XYX`$, which case is treated above.
These calculations verify that $`L(c𝒜_v)`$ can be identified with $`L(\overline{T}_X(G))`$ as described above, with the same rank function. ∎
###### Remark 2.5
The same calculations as in case 1 above can be used to show that the converse of Theorem 2.4 also holds. That is, $`X`$ is modular if the lattice $`L(c𝒜_{\overline{v}})`$ is constant over $`M(𝒜_X)`$. This will be used to identify modular flats in the examples of Section 5.∎
Let $`M(𝒜)=^{\mathrm{}}𝒜`$ and $`M(𝒜_X)=(^{\mathrm{}}/X)𝒜_X`$. Note that $`\pi `$ maps $`M(𝒜)`$ onto $`M(𝒜_X)`$.
###### Corollary 2.6
The fibers of $`\pi |_{M(𝒜)}:M(𝒜)M(𝒜_X)`$ are diffeomorphic.
###### Proof.
The fiber of $`\pi |_{M(𝒜)}`$ over $`\overline{v}`$ is the complement of the arrangement $`𝒜_{\overline{v}}`$ in $`\pi ^1(\overline{v})^{\mathrm{}r(X)}`$. Since the base $`M(𝒜_X)`$ is path-connected, Theorem 2.4 implies that the arrangements $`c𝒜_v`$ are lattice-isotopic. Then the assertion follow from . ∎
###### Remark 2.7
Theorem 2.4 was essentially proved by Terao in . Our result explicitly identifies the lattice. In case $`X`$ is a copoint, Corollary 2.6 follows without using Randell’s lattice isotopy theorem, which had not been discovered at the time of Terao’s work. In fact Corollary 2.6 and the fibration result Corollary 3.2 of the next section confirm the suggestion stated after Proposition 2.12 of . The proof of Corollary 3.2 uses the stratification technique first introduced to arrangement theory by Randell in his proof of the isotopy theorem.
## 3 Modular flats and fibrations
The arrangement $`𝒜`$ defines a stratification $`𝒮`$ of $`^{\mathrm{}}`$:
$$^{\mathrm{}}=\{S_Y|YL\}$$
whose strata $`S_Y`$ are given by
$$S_Y=Y\underset{Z>Y}{}Z.$$
Thus $`S_Y`$ is a connected dense open subset of the linear space $`Y`$. In particular, $`S_Y`$ is a smooth submanifold of $`^{\mathrm{}}`$. Note that the closed stratum $`\overline{S_Y}`$ is equal to $`Y`$. Also $`S_Y\overline{S_Z}\mathrm{}`$ if and only if $`S_Y\overline{S_Z}`$ if and only if $`YZ`$.
This stratification satisfies Whitney’s conditions (a) and (b) . Indeed these conditions involve tangent and secant lines, and tangent spaces to strata, which are trivial to verify because $`S_Y`$, as an open subset of the linear space $`Y`$, has tangent space at any point equal to $`Y`$.
Let $`XL`$ be a modular element of rank $`p`$. We may identify $`^{\mathrm{}}/X`$ with $`^p`$. Let $`\pi :^{\mathrm{}}^p`$ be the natural projection. The arrangement $`𝒜_X=\{H𝒜|HX\}`$, considered as an arrangement in $`^{\mathrm{}}/X^p`$, determines a stratification of $`^p`$ as above. Elements of $`L(𝒜_X)`$ have the form $`\pi Y=X+Y/X`$ for $`YL(𝒜)`$. Referring to Lemma 2.3, one sees that the preimage of a stratum is a union of strata, that is, that $`\pi `$ is a map of stratified spaces, precisely when $`X`$ is modular. Since $`\pi `$ is a linear surjection, it restricts to a submersion on each stratum. In order to apply the Thom Isotopy Lemma, it is necessary to extend $`\pi `$ to a proper map of stratified spaces. This step was carried out by L. Paris .
###### Theorem 3.1
There exists a stratified space $`P_X`$ containing $`^{\mathrm{}}`$ as an open dense subset, and an extension of $`\pi `$ to a proper stratified map $`\widehat{\pi }:P_X^p`$.∎
The space $`P_X`$ is obtained by compactifying the fibers of $`\pi `$, i.e., the parallel translates of $`X`$, via projective completion, so that $`P_X`$ is diffeomorphic to $`(^q)\times ^p`$, where $`q=\mathrm{}p=dim(X)`$. This can be viewed as a parametrized version of R. Randell’s construction in his proof of the lattice isotopy theorem . The stratification of $`^{\mathrm{}}`$ is extended to a stratification of $`P_X`$ by adjoining closed strata formed by intersecting the closures of the $`S_Y`$ in $`P_X`$ with $`((^q)^q)\times ^p`$. These new strata have the form $`S_Y^{\mathrm{}}\times ^p`$, for $`YX`$, where $`S_Y^{\mathrm{}}=(S_YX)((^q)^q)`$. The map $`\widehat{\pi }`$ is projection on the second factor.
Let $`M(𝒜)`$ and $`M(𝒜_X)`$ denote the complements of $`𝒜`$ and $`𝒜_X`$ in $`^{\mathrm{}}`$ and $`^p`$ respectively.
###### Corollary 3.2
The map $`\pi |_{M(𝒜)}:M(𝒜)M(𝒜_X)`$ is a fiber bundle projection.
###### Proof.
The complement $`M(𝒜)`$ coincides with the open stratum $`S_{0_L}`$ of $`^{\mathrm{}}P_X`$. So Theorem 3.1 implies that the restriction of $`\widehat{\pi }`$ to $`M(𝒜)`$ is a fiber bundle projection, by the Thom Isotopy Lemma . ∎
We proceed to generalize the properties of strictly linear fibrations , where $`X`$ is a modular copoint, to general modular fibrations. Henceforth let $`X`$ be a modular flat of $`L(𝒜)`$, and let us denote the bundle projection $`\pi |_{M(𝒜)}`$ by $`\pi _X`$.
We say $`𝒜`$ is a $`K(\pi ,1)`$ arrangement if $`M(𝒜)`$ is an aspherical space.
###### Corollary 3.3
If $`𝒜_X`$ and the coned fiber arrangement $`c𝒜_{\overline{v}}`$ are $`K(\pi ,1)`$ arrangements, then $`𝒜`$ is a $`K(\pi ,1)`$ arrangement.
###### Proof.
This follows immediately from the long exact homotopy sequence of the fibration $`\pi _X`$. ∎
###### Remark 3.4
In case $`X`$ is a modular copoint, the monodromy of $`\pi _X`$ induces a homomorphism from $`\pi _1(M(𝒜_X))`$ to $`P_n`$, the pure braid group on $`n=|𝒜𝒜_X|`$ strands, which we call the braid monodromy homomorphism after its similarity to the Moishezon construction. See . For a modular flat $`X`$ of arbitrary rank, the pure braid group is replaced by the fundamental group of a certain subvariety of the Grassmanian, a matroid stratum defined as follows. If $`P𝒢_{\mathrm{}}(^n)`$ is a point of the Grassmannian of $`\mathrm{}`$-planes in $`^n`$, then $`P`$ determines a vector configuration in $`^{\mathrm{}}`$, unique up to linear change of coordinates, obtained by projecting the standard basis vectors of $`^n`$ onto $`P`$ . Let $`G_P`$ denote the linear matroid realized by this configuration; $`G_P`$ is independent of the choice of basis in $`P`$. The matroid stratum of an arbitrary matroid $`G`$ is the subset $`\mathrm{\Gamma }(G)`$ of $`𝒢_{\mathrm{}}(^n)`$ given by
$$\mathrm{\Gamma }(G)=\{P𝒢_{\mathrm{}}(^n)|G_P=G\}.$$
An ordered arrangement $`𝒜=\{H_1,\mathrm{},H_n\}`$ of rank $`\mathrm{}`$, with specified defining forms $`\{\varphi _1,\mathrm{},\varphi _n\}`$, determines a point $`P𝒢_{\mathrm{}}(^n)`$ given by the image of $`(\varphi _1,\mathrm{},\varphi _n):^{\mathrm{}}^n.`$ The original arrangement $`𝒜`$ is isomorphic to the arrangement in $`P`$ formed by the intersections of $`P`$ with the coordinate hyperplanes in $`^n`$, and the point $`P`$ lies in $`\mathrm{\Gamma }(G(𝒜))`$. See .
The monodromy of the stratified map $`\pi `$ induces a homomorphism
$$\pi _1(M(𝒜_X))\pi _1(\mathrm{\Gamma }(\overline{T}_X(G))).$$
Indeed, a path $`\{\overline{v}_t\}_{t[0,1]}`$ in the base space $`M(𝒜_X)`$ determines a one-parameter family of (coned) fiber arrangements $`c𝒜_{\overline{v}_t}`$, equipped with ordered sets of defining forms inherited from a fixed set of defining forms for $`𝒜`$. By Theorem 2.4 and the construction above, this defines a path in the matroid stratum $`\mathrm{\Gamma }(\overline{T}_X(G))`$. From this one easily obtains the monodromy homomorphism described above.
This construction does indeed generalize the corank one case. For in this case $`\overline{T}_X(G)`$ is a uniform matroid of rank two, $`\mathrm{\Gamma }(\overline{T}_X(G))`$ is configuration space, and $`\pi _1(\mathrm{\Gamma }(\overline{T}_X(G)))`$ is the pure braid group.∎
###### Theorem 3.5
The monodromy action of $`\pi _1(M(𝒜_X))`$ on the fiber $`M(𝒜_{\overline{v}})`$ is cohomologically trivial.
###### Proof.
Since the fiber $`M(𝒜_{\overline{v}})`$ is the complement of an arrangement, the cohomology of $`M(𝒜_{\overline{v}})`$ is free abelian, and is generated by $`H^1(M(𝒜_{\overline{v}}))`$. First of all we argue that the monodromy action on $`H^1(M(𝒜_{\overline{v}}))`$ is trivial, by the same reasoning as in the corank-one case . The group $`H^1(M(𝒜_{\overline{v}}))`$ has a free basis consisting of elements dual to the hyperplanes of $`𝒜_{\overline{v}}`$. Using this basis, it is clear that elements of $`H^1(M(𝒜_{\overline{v}}))`$ are uniquely determined by their linking numbers with the hyperplanes of $`𝒜_{\overline{v}}`$. By naturality, these linking numbers agree with linking numbers in $`^{\mathrm{}}`$ with the hyperplanes of $`𝒜𝒜_X`$. Since these linking numbers take values in a discrete space, and vary continuously, they remain locally constant under translation of the fiber, and thus are globally constant under translation around a loop in the base. This proves triviality in degree one. Since $`H^{}(M(𝒜_{\overline{v}}))`$ is generated by $`H^1(M(𝒜_{\overline{v}}))`$, and the monodromy action respects cup products, it follows that the monodromy acts trivially on $`H^{}(M(𝒜_{\overline{v}}))`$. ∎
A rational $`K(\pi ,1)`$ arrangement is an arrangement whose complement has aspherical rational completion. See for the precise definition and basic properties. We point out that this property seems to bear little relationship to the notion of $`K(\pi ,1)`$ arrangement; the terminology arises naturally in the context of simply-connected spaces.
###### Corollary 3.6
If $`𝒜_X`$ and $`c𝒜_{\overline{v}}`$ are rational $`K(\pi ,1)`$ arrangements, then $`𝒜`$ is a rational $`K(\pi ,1)`$ arrangement.
###### Proof.
The argument is the same as in the corank-one case . Because the monodromy action is trivial, hence nilpotent, on the cohomology of the fiber, the map $`\pi _X`$ induces a fibration of the rational completion of $`M(𝒜)`$ over that of $`M(𝒜_X)`$, with fiber the rational completion of $`M(𝒜_{\overline{v}})`$. Since $`M(c𝒜_{\overline{v}})^{}\times M(𝒜_{\overline{v}})`$, the hypothesis implies that the rational completion of $`M(𝒜_{\overline{v}})`$ is aspherical. The assertion then follows from the homotopy sequence of this fibration. ∎
At this point the only known examples of rational $`K(\pi ,1)`$ arrangements are supersolvable. If $`𝒜_X`$ and $`c𝒜_{\overline{v}}`$ are supersolvable, then $`𝒜`$ is also supersolvable . So the preceding corollary does not provide new examples of rational $`K(\pi ,1)`$ arrangements.
The Poincaré series of a topological space $`M`$ is
$$P(M,t)=\underset{n0}{}dim_{}H^n(M,).$$
For a complex arrangement $`𝒜`$, a famous result of Orlik and Solomon relates the Poincaré series $`P(M(𝒜),t)`$ to the characteristic polynomial $`\chi (G(𝒜),t)`$ of the underlying matroid $`G(𝒜)`$. Specifically,
$$P(M(𝒜),t)=t^r\chi (G(𝒜),t^1),$$
where $`r`$ is the rank of $`G(𝒜)`$.
For a modular flat $`X`$, R. Stanley proved in that the characteristic polynomial of the $`G(𝒜_X)`$ divides that of $`G(𝒜)`$ over the integers. In , T. Brylawski identified the quotient as the characteristic polynomial of the complete principal truncation $`\overline{T}_X(G)`$, divided by $`(t1)`$. The decone operation on arrangements has the effect on Poincaré polynomials of dividing by $`(1+t)`$. Using Theorem 2.4 and the identity relating the characteristic polynomial of $`G(𝒜)`$ to the Poincaré polynomial of $`M(𝒜)`$, we may restate the Stanley and Brylawski results as follows.
###### Theorem 3.7
If $`X`$ is a modular flat of $`G`$, then
$$P(M(𝒜),t)=P(M(𝒜_X),t)P(M(𝒜_{\overline{v}}),t).$$
###### Corollary 3.8
The Leray-Serre spectral sequence of $`\pi _X`$ satisfies
$$E_2^{p,q}H^p(M(𝒜_X))H^q(M(𝒜_{\overline{v}})),$$
and degenerates at the $`E_2`$ term.∎
###### Proof.
The first assertion follows from the triviality of the monodromy action established in Theorem 3.5. The second is a consequence of the factorization identity among the Poincaré series. Indeed, according to \[17, Theorem 11.3\], the formula of Theorem 3.7 holds for a general spectral sequence $`E`$, with a correction term that vanishes precisely when the differential of $`E_2`$ is trivial. ∎
###### Corollary 3.9
The cohomology $`H^{}(M(𝒜))`$ is isomorphic as a $``$-module to the tensor product $`H^{}(M(𝒜_X))H^{}(M(𝒜_{\overline{v}}))`$.∎
###### Remark 3.10
In Terao established the tensor product factorization of Corollary 3.9 in terms of Orlik-Solomon algebras, using a direct combinatorial argument. This approach yields an alternate proof of the Stanley factorization theorem. The degeneracy of the spectral sequence in case $`X`$ is a modular copoint is given a direct proof in , providing a topological proof of Terao’s result in this case. The proof in uses the fact that the fiber $`M(𝒜_{\overline{v}})`$ has nonvanishing cohomology in only two different degrees, so that the spectral sequence results in a “Gysin-like” long exact sequence. The other ingredient is the construction of a section of the bundle map $`\pi _X`$. In case $`X`$ is a modular flat of arbitrary rank, a section of $`\pi _X`$ is constructed by L. Paris in . But we see no analogue of the Gysin long exact sequence in the general case, and do not have a topological proof, independent of the Stanley and Brylawski results, of the second part of Corollary 3.8. Nevertheless, the bundle map $`\pi _X`$ is seen to be a topological realization of the combinatorial and algebraic factorizations arising from a modular flat.∎
Motivated by the fact that supersolvable arrangements are inductively free , we include with this compendium of generalizations the following conjecture.
###### Conjecture 3.11
If $`X`$ is a modular flat and both $`𝒜_X`$ and $`c𝒜_{\overline{v}}`$ are free arrangements, then $`𝒜`$ is a free arrangement.
## 4 Parallel connections
Let $`G_1`$ and $`G_2`$ be matroids on ground sets $`E_1`$ and $`E_2`$. Suppose $`E_1E_2=F`$ is a flat of both $`G_1`$ and $`G_2`$, and is modular in $`G_1`$. The generalized parallel connection of $`G_1`$ and $`G_2`$ along $`F`$ is the matroid $`P_F(G_1,G_2)`$ on the ground set $`E_1E_2`$ whose flats are those sets $`YE_1E_2`$ for which $`YE_i`$ is a flat of $`G_i`$ for $`i=1,2`$. The modularity condition is necessary for this definition to make sense. That is, this collection of flats will form a geometric lattice for general $`G_2`$ if and only if $`X`$ is modular in $`G_1`$. Modularity of $`F`$ in $`G_1`$ implies that $`G_2`$ is modular in $`P_F(G_1,G_2)`$. See \[36, Section 7.6\] and for details about this construction.
The rank of a flat $`Y`$ of $`P_F(G_1,G_2)`$ is given by
$$r(Y)=r_1(YE_1)+r_2(YE_2)r_1(YF),$$
where $`r_i`$ is the rank function of $`G_i,i=1,2`$. In particular, the rank of $`P_F(G_1,G_2)`$ is equal to $`r(G_1)+r(G_2)r(F).`$ The rank formula indicates that $`P_F(G_1,G_2)`$ is the “free” sum of $`G_1`$ and $`G_2`$ amalgamated along their common flat $`F`$. Indeed, $`P_F(G_1,G_2)`$ is a pushout of the inclusion maps $`FG_i,i=1,2`$ in the category of matroids and injective strong maps.
In case $`F`$ is a point, automatically modular in $`G_1`$, the matroid $`P_F(G_1,G_2)`$ is called a parallel connection of $`G_1`$ and $`G_2`$, studied in connection with complex hyperplane arrangements in .
Now suppose $`𝒜_1`$ and $`𝒜_2`$ are hyperplane arrangements realizing $`G_1`$ and $`G_2`$ in $`^r`$ and $`^s`$ respectively. Then there is an arrangement $`𝒜`$ realizing $`P_F(G_1,G_2)`$, provided there is a linear isomorphism between the subarrangements of $`𝒜_1`$ and $`𝒜_2`$ corresponding to the common flat $`F`$. To carry out the construction, let us be more precise about the realizations $`𝒜_1`$ and $`𝒜_2`$.
Suppose the flat $`F`$ has rank $`p`$ (in both $`G_1`$ and $`G_2`$). Let $`X_1`$ denote the corresponding element of intersection lattice $`L(𝒜_1)`$. Thus $`X_1`$ is a linear subspace of $`^r`$, and we may identify $`(𝒜_1)_{X_1}`$ with $`F`$. We may assume $`X_1=^{rp}\times \{0\}^r`$. Then the defining equations of the hyperplanes in $`(𝒜_1)_{X_1}𝒜_1`$ involve only the last $`p`$ coordinates in $`^r`$. Assume that the same defining forms, expressed in terms of the first $`p`$ coordinates of $`^s`$, give the defining equations for hyperplanes of $`(𝒜_2)_{X_2}𝒜_2`$, where $`X_2L(𝒜_2)`$ corresponds to the flat $`F`$ of $`G_2`$. Then we may define an arrangement $`𝒜`$ in $`^{\mathrm{}}`$, with $`\mathrm{}=r+sp`$, as follows. Identify $`^{\mathrm{}}`$ with $`^{rp}\times ^p\times ^{sp}`$. By pulling back the defining equations via projection of coordinates, the arrangements $`𝒜_1`$ and $`𝒜_2`$ naturally embed in $`^r\times ^{sp}=^{\mathrm{}}`$ and $`^{rp}\times ^s=^{\mathrm{}}`$ respectively. Then let $`𝒜`$ be the union of $`𝒜_1(𝒜_1)_{X_1}`$ and $`𝒜_2`$ in $`^{\mathrm{}}`$.
###### Theorem 4.1
\[36, Prop. 7.6.11\] The arrangement $`𝒜`$ is a realization of the generalized parallel connection $`P_F(G_1,G_2)`$. ∎
Let $`XL(𝒜)`$ correspond to the flat $`F`$ of $`P_F(G_1,G_2)`$. By modularity of $`F`$ in $`G_1`$ and of $`G_2`$ in $`P_F(G_1,G_2)`$, the results of Section 3 yield bundle maps $`M(𝒜_1)M((𝒜_1)_{X_1})`$ and $`M(𝒜)M(𝒜_2)`$. We consider $`(𝒜_1)_{X_1}`$ to be an arrangement in $`^r/X_1^s/X_2^p`$. Then there is a projection $`M(𝒜_2)M((𝒜_1)_{X_1})`$. This projection is just the inclusion of $`M(𝒜_2)`$ into the complement of the subarrangement $`\{H𝒜_2|HX_2\}`$ of $`𝒜_2`$, followed by a homotopy equivalence.
###### Theorem 4.2
The fiber bundle $`M(𝒜)M(𝒜_2)`$ is the pullback of the bundle $`M(𝒜_1)M((𝒜_1)_{X_1})`$ along the projection $`M(𝒜_2)M((𝒜_1)_{X_1})`$. That is,
$$\begin{array}{ccc}M(𝒜)& & M(𝒜_1)\\ & & & & \\ M(𝒜_2)& & M((𝒜_1)_{X_1})\end{array}$$
is a pullback diagram.
###### Proof.
For these special realizations, the bundle map $`M(𝒜_1)M((𝒜_1)_{X_1})`$ is the restriction of the projection $`\pi _1:^r^p`$ onto the last $`p`$ coordinates. Similarly, the map $`M(𝒜)M(𝒜_2)`$ is the restriction of the projection $`\pi :^{\mathrm{}}^s`$ onto the last $`s`$ coordinates. The map $`M(𝒜_2)M((𝒜_1)_{X_1})`$ can be identified with the restriction of the projection $`\pi _2:^sC^s/X_2^p`$ onto the first $`p`$ coordinates.
By definition, the total space of the pullback of $`M(𝒜_1)M((𝒜_1)_{X_1})`$ along the projection $`M(𝒜_2)M(𝒜_X)`$ is the set of pairs $`(x,v)^r\times ^s`$ such that $`xM(𝒜_1),vM(𝒜_2)`$, and $`\pi _1(x)=\pi _2(v)`$ in $`M((𝒜_1)_{X_1})C^p`$. But this means that the last $`p`$ components of $`x`$ match the first $`p`$ components of $`v`$. Then each such $`(x,v)`$ corresponds to a unique point of $`^{r+sp}=^{\mathrm{}}`$ which, by the first two conditions, lies in $`M(𝒜)`$. Under this identification, the projection $`(x,v)v`$ coincides with $`\pi `$. This identifies $`\pi |_{M(𝒜)}`$ with the pullback of $`\pi _1|_{M(𝒜_1)}`$, as claimed. ∎
###### Corollary 4.3
If $`𝒜`$ is a realization of the parallel connection of $`𝒜_1`$ and $`𝒜_2`$, then $`M(𝒜)`$ is a trivial bundle over $`M(𝒜_2)`$ with fiber $`M(d𝒜_1)`$. In particular, $`M(𝒜)M(d𝒜_1)\times M(𝒜_2)`$.
###### Proof.
In case $`X`$ is a point, then $`X`$ is modular in $`𝒜_1`$, and the modular fibration $`M(𝒜_1)M(𝒜_X)=^{}`$ is a trivial bundle with fiber $`d𝒜_1`$, by \[23, Proposition 5.1\]. The pullback of a trivial bundle is trivial. ∎
This corollary clarifies the main construction of , which essentially established the diffeomorphism noted above. This argument shows in an alternate way that the diffeomorphisms among arrangements with different underlying matroids, constructed in , are all consequences of the triviality of the restriction of the Hopf bundle.
## 5 Examples
Corollary 3.3 of Section 3 can be used to identify $`K(\pi ,1)`$ arrangements of high rank, at least when the base arrangement $`𝒜_X`$ and (coned) fiber arrangement $`c𝒜_{\overline{v}}`$ are tractable. This will be the case, for instance, when $`X`$ is a modular coline, for then $`c𝒜_{\overline{v}}`$ will have rank three. In this section we present new families of examples of $`K(\pi ,1)`$ arrangements. Our results give some support for the conjecture , which was based primarily on rank-three phenomena, that factored arrangements of arbitrary rank are $`K(\pi ,1)`$. We also exhibit an interesting example with two different modular colines, allowing us to conclude the nontrivial result that one of the fiber arrangements is not $`K(\pi ,1)`$.
Let $`_{\mathrm{}}`$ denote the arrangement of reflecting hyperplanes in the Weyl group of type $`B_{\mathrm{}}`$. Thus $`_{\mathrm{}}`$ consists of the hyperplanes
$$H_{ij}=\{x^{\mathrm{}}|x_i=x_j\},\text{for}1i<j\mathrm{},$$
$$\overline{H}_{ij}=\{x^{\mathrm{}}|x_i=x_j\}\text{for}1i<j\mathrm{},\text{and}$$
$$H_i=\{x^{\mathrm{}}|x_i=0\},\text{for}1i\mathrm{}.$$
Let $`𝒜_\mathrm{}1`$ denote the braid arrangement, consisting of the hyperplanes $`H_{ij}`$ above, for $`1i<j\mathrm{}`$.
In , P. Edelman and V. Reiner used graphs to parametrize subarrangements of $`_{\mathrm{}}`$ containing $`𝒜_\mathrm{}1`$, developing a calculus for combinatorial invariants of the arrangements in terms of the graphs. We find among these arrangements those which are not supersolvable, but have modular colines, for which the fiber arrangements are demonstrably $`K(\pi ,1)`$. These examples coincide in large part with the arrangements between $`𝒜_\mathrm{}1`$ and $`_{\mathrm{}}`$ which are factored, classified by D. Bailey in .
Let $`\mathrm{\Gamma }`$ be a graph with vertex set $`\{1,\mathrm{},\mathrm{}\}`$, possibly with loops, but without multiple edges. Let $`\mathrm{edge}(\mathrm{\Gamma })`$ and $`\mathrm{loop}(\mathrm{\Gamma })`$ denote the sets of edges and loops of $`\mathrm{\Gamma }`$, respectively. Let $`𝒜_\mathrm{\Gamma }`$ be the arrangement defined by
$$𝒜_\mathrm{\Gamma }=𝒜_\mathrm{}1\{\overline{H}_{ij}|ij\mathrm{edge}(\mathrm{\Gamma })\}\{H_i|i\mathrm{loop}(\mathrm{\Gamma })\}.$$
The following results are proved in . The notion of free arrangement plays little role in what follows; see for a precise definition. A graph is threshold if it is built up by successively adjoining isolated and/or cone vertices, the latter being vertices which are adjacent to all preceding vertices.
###### Theorem 5.1
The arrangement $`𝒜_\mathrm{\Gamma }`$ is free if and only if
1. $`\mathrm{\Gamma }`$ is a threshold graph, and
2. $`i\mathrm{loop}(\mathrm{\Gamma })`$ and $`\mathrm{deg}(j)>\mathrm{deg}(i)`$ implies $`j\mathrm{loop}(\mathrm{\Gamma })`$. ∎
An edge $`ijedge(\mathrm{\Gamma })`$ is loopless if neither $`i`$ nor $`j`$ lies in $`\mathrm{loop}(\mathrm{\Gamma })`$.
###### Theorem 5.2
Suppose $`𝒜_\mathrm{\Gamma }`$ is free and $`\mathrm{\Gamma }`$ has no loopless edges. Then $`𝒜_\mathrm{\Gamma }`$ is supersolvable.∎
There are two families of exceptional graphs $`\mathrm{\Gamma }`$ with loopless edges such that $`𝒜_\mathrm{\Gamma }`$ is supersolvable; see . Of course, any such arrangement $`𝒜_\mathrm{\Gamma }`$ is $`K(\pi ,1)`$.
Roughly speaking, an arrangement is factored if the cohomology of the complement is isomorphic as a $``$-module to the the tensor product of algebras with trivial multiplication generated by sets of hyperplanes of $`𝒜`$. For instance, Corollary 3.9 implies that supersolvable arrangements are factored. Such factorizations correspond to partitions of $`𝒜`$, properties of which are analyzed in . D. Bailey identified those arrangements $`𝒜_\mathrm{\Gamma }`$ which are factored.
###### Theorem 5.3
Suppose $`𝒜_\mathrm{\Gamma }`$ is free and $`\mathrm{\Gamma }`$ has at most one loopless edge. Then $`𝒜_\mathrm{\Gamma }`$ is a factored arrangement. ∎
Again there are two families of exceptional graphs with more than one loopless edge for which $`𝒜_\mathrm{\Gamma }`$ is factored .
We establish criteria for $`𝒜_\mathrm{\Gamma }`$ to have modular copoints or colines determined by coordinate subspaces $`X`$. The assertions below are easy to prove using Remark 2.5, by showing that the lattices of the coned fiber arrangements remain constant over $`M(𝒜_X)`$. (For an example, see the proof of Theorem 5.6.) Let $`\mathrm{\Gamma }^{}`$ and $`\mathrm{\Gamma }^{\prime \prime }`$ be the vertex-induced subgraphs of $`\mathrm{\Gamma }`$ on vertices $`\{2,\mathrm{},\mathrm{}\}`$ and $`\{3,\mathrm{},\mathrm{}\}`$ respectively. Then $`𝒜_\mathrm{\Gamma }^{}`$ and $`𝒜_{\mathrm{\Gamma }^{\prime \prime }}`$ are flats of $`G(𝒜_\mathrm{\Gamma })`$ of corank one and two corresponding to the linear subspaces $`x_2=\mathrm{}=x_{\mathrm{}}=0`$ and $`x_3=\mathrm{}=x_{\mathrm{}}=0`$ respectively.
###### Lemma 5.4
The flat $`𝒜_\mathrm{\Gamma }^{}`$ is a modular copoint of $`G(𝒜_\mathrm{\Gamma })`$ if and only if
1. $`1\mathrm{loop}(\mathrm{\Gamma })`$ implies $`\mathrm{loop}(\mathrm{\Gamma })=\{1,\mathrm{},\mathrm{}\}`$, and
2. $`1j\mathrm{edge}(\mathrm{\Gamma })`$ implies $`j\mathrm{loop}(\mathrm{\Gamma })`$ and $`j`$ is adjacent
to every vertex of $`\mathrm{\Gamma }`$. ∎
In particular, an isolated vertex of $`\mathrm{\Gamma }`$ corresponds to a modular copoint.
###### Lemma 5.5
The flat $`𝒜_{\mathrm{\Gamma }^{\prime \prime }}`$ is a modular coline of $`G(𝒜_\mathrm{\Gamma })`$ if and only if
1. $`i\mathrm{loop}(\mathrm{\Gamma })`$, for $`i=1`$ or $`2`$ implies $`\mathrm{loop}(\mathrm{\Gamma })\{3,\mathrm{},\mathrm{}\}`$,
2. $`ij\mathrm{edge}(\mathrm{\Gamma })`$ for $`i=1`$ or $`2`$ implies $`j\mathrm{loop}(\mathrm{\Gamma })`$ and $`j`$ is adjacent to every vertex $`k`$ for $`k3`$, and
3. $`12\mathrm{edge}(\mathrm{\Gamma })`$ implies $`\mathrm{loop}(\mathrm{\Gamma })\{3,\mathrm{},\mathrm{}\}`$. ∎
The modular fibration corresponding to a modular copoint has as coned fiber arrangement $`c𝒜_{\overline{v}}`$ a central arrangement of rank two, which is $`K(\pi ,1)`$, so $`𝒜_\mathrm{\Gamma }`$ is $`K(\pi ,1)`$ if and only if $`𝒜_\mathrm{\Gamma }^{}`$ is $`K(\pi ,1)`$. This is a “strictly linear fibration” as studied in .
Factored arrangements of rank three were shown to be $`K(\pi ,1)`$ in . Supersolvable arrangements of arbitrary rank are factored, and are $`K(\pi ,1)`$. In we conjecture that factored arrangements of arbitrary rank are $`K(\pi ,1)`$. The next result provides more support for this conjecture. The arrangements of this theorem are not supersolvable, by Theorem 5.2, but are factored, by Theorem 5.3. In fact, by an argument in , the examples described below are the only factored non-supersolvable arrangements $`𝒜_\mathrm{\Gamma }`$ which have only two non-loop vertices.
###### Theorem 5.6
Suppose $`\mathrm{\Gamma }`$ is the complete graph on $`\mathrm{}`$ vertices, and $`\mathrm{loop}(\mathrm{\Gamma })=\{3,\mathrm{},\mathrm{}\}`$. Then $`𝒜_\mathrm{\Gamma }`$ is a $`K(\pi ,1)`$ arrangement.
###### Proof.
Let $`𝒜=𝒜_\mathrm{\Gamma }`$. Let $`\mathrm{\Gamma }^{\prime \prime }`$ denote the vertex-induced subgraph of $`\mathrm{\Gamma }`$ on vertices $`\{3,\mathrm{},\mathrm{}\}`$. Then $`𝒜_{\mathrm{\Gamma }^{\prime \prime }}`$ is a modular coline of $`G(𝒜)`$ by Lemma 5.5, and $`𝒜_{\mathrm{\Gamma }^{\prime \prime }}`$ is supersolvable, hence $`K(\pi ,1)`$, by Theorem 5.2. Let $`X=𝒜_{\mathrm{\Gamma }^{\prime \prime }}`$, and let $`\overline{v}=(v_3,\mathrm{},v_{\mathrm{}})M(𝒜_X).`$ Then $`v_i\pm v_j`$ for $`3i<j\mathrm{}`$, and $`v_j0`$ for $`3j\mathrm{}`$. The fiber arrangement $`𝒜_{\overline{v}}`$ is the affine arrangement in $`^2`$ consisting of the lines
$$x_1=\pm x_2,x_1=\pm v_j,\text{and}x_2=\pm v_j$$
for $`3j\mathrm{}.`$
The coned fiber arrangement $`c𝒜_{\overline{v}}`$ pictured in Figure 1 is a $`K(\pi ,1)`$ arrangement. Indeed $`c𝒜_{\overline{v}}`$ is precisely Example 3.13 of . Alternatively, $`c𝒜_{\overline{v}}`$ is a factored arrangement of rank three, hence $`K(\pi ,1)`$ by . We conclude by Corollary 3.3 that $`𝒜`$ is $`K(\pi ,1)`$. ∎
###### Remark 5.7
The conclusion of the theorem also holds if $`\mathrm{loop}(\mathrm{\Gamma })=\mathrm{}`$, for then $`𝒜_\mathrm{\Gamma }`$ is a Coxeter arrangement of type $`D_{\mathrm{}}`$, which is simplicial.∎
One can use Theorem 5.6 and Lemma 5.4 to build other examples of non-supersolvable (and non-simplicial) $`K(\pi ,1)`$ arrangements of high rank, by successively adding vertices satisfying 5.4 to the graphs of Theorem 5.6. See Figure 2. The existence of loops in $`\mathrm{\Gamma }`$ is essential: the same construction with the loopless complete graph (the $`D_{\mathrm{}}`$ arrangement) allows only the addition of isolated vertices.
There is one rank-three arrangement $`𝒜_\mathrm{\Gamma }`$, a realization of the non-Fano matroid, which is not supersolvable, but is simplicial, hence $`K(\pi ,1)`$. This arrangement can also be used with Theorem 5.4 to construct non-supersolvable $`K(\pi ,1)`$ arrangements. This construction is illustrated in Figure 3.
###### Example 5.8
We exhibit in Figures 2 and 3 some other graphs $`\mathrm{\Gamma }`$ for which $`𝒜_\mathrm{\Gamma }`$ is $`K(\pi ,1)`$, using the constructions of the preceding paragraphs. These arrangements are factored, as is every arrangement arising from these constructions, by Theorem 5.3. ∎
We close with another interesting example from the class of “$`A`$-$`B`$ arrangements.”
###### Example 5.9
Let $`\mathrm{\Gamma }`$ be the graph with vertex set $`\{1,2,3,4\}`$, edges $`12,13,`$ and $`24`$, and loops at vertices $`1,2,\text{and}3`$. Let $`\mathrm{\Gamma }_0^{\prime \prime }`$ and $`\mathrm{\Gamma }_1^{\prime \prime }`$ be the vertex-induced subgraphs of $`\mathrm{\Gamma }`$ with vertex sets $`\{1,2\}`$ and $`\{1,3\}`$ respectively. Then both $`𝒜_{\mathrm{\Gamma }_0^{\prime \prime }}`$ and $`𝒜_{\mathrm{\Gamma }_1^{\prime \prime }}`$ are modular flats of $`G(𝒜_\mathrm{\Gamma })`$. The respective fiber arrangements are illustrated in Figure 4.
The arrangement on the right, the (coned) fiber arrangement of $`M(𝒜_\mathrm{\Gamma })M(𝒜_{\mathrm{\Gamma }_1^{\prime \prime }})`$, is not $`K(\pi ,1)`$ because it has a “simple triangle” \[14, Corollary 3.3\]. Now $`𝒜_{\mathrm{\Gamma }_1^{\prime \prime }}`$ is a $`K(\pi ,1)`$ arrangement, being a central arrangement of rank two. It follows that $`M(𝒜_\mathrm{\Gamma })`$ is not aspherical. But $`𝒜_{\mathrm{\Gamma }_0^{}}`$ is also $`K(\pi ,1)`$. We conclude that the (coned) fiber arrangement of $`M(𝒜_\mathrm{\Gamma })M(𝒜_{\mathrm{\Gamma }_0^{}})`$, shown on the left, cannot be $`K(\pi ,1)`$. This is the only argument we know of to show this arrangement is not $`K(\pi ,1)`$. ∎
The research presented here, and our general interest in $`K(\pi ,1)`$ arrangements, was motivated in part by a suggestion of G. Ziegler of a straightforward construction of rank-three arrangements with the same underlying matroid but homotopy inequivalent complements. The argument avoids fundamental group computations, but relies on the existence of high-rank $`K(\pi ,1)`$ arrangements with certain properties, whose existence has not yet been shown. Here is the precise problem, to which the methods of this paper may apply.
###### Problem 5.10
(Ziegler) Find $`K(\pi ,1)`$ arrangements whose matroids have the same flats of ranks one and two but have different characteristic polynomials.
###### Acknowledgements
We are grateful to Joseph Kung, Luis Paris, Hiroaki Terao and Vic Reiner and David Bailey, for helpful conversations and correspondence, and to Terence Blows, who organized the REU program at Northern Arizona University, where this project was begun.
Department of Mathematics and Statistics
Northern Arizona University
Flagstaff, AZ 86011-5717
michael.falk@nau.edu
Department of Mathematics
Harvard University
Cambridge, MA
Proudf@fas.harvard.edu |
warning/0002/astro-ph0002419.html | ar5iv | text | # On The Perturbations Of Viscous Rotating Newtonian Fluids
## 1 INTRODUCTION
The problem of perturbations of rotating fluids is of considerable interest in astrophysics. These perturbations can give us information not only about stars and accretion disks (e.g., Friedman & Ipser 1992; Nowak & Wagoner 1992; Kato et al. 1998), but also about the properties of the compact object at the center of the latter (e.g., Perez et al. 1997).
Traditionally, this problem has been approached in two different ways. The classical way has been to use a Lagrangian displacement formalism (Lynden-Bell & Ostriker 1967). An alternative method uses two scalar potentials (proportional to the perturbations of the pressure and gravitational potential) within the framework of Eulerian fluctuations (Poincaré 1885; Ipser & Managan 1985; Managan 1985). In the spirit of the latter approach, Ipser & Lindblom (1991b, hereafter IL; 1992) have developed an elegant formalism to analyze adiabatic pulsations of perfect Newtonian (as well as relativistic) fluids that are stationary, axisymmetric, and purely rotating (i.e., have no meridional component of velocity). They look for normal-mode solutions and reduce the problem to two coupled second-order partial differential equations for the two potentials.
The inclusion of viscosity in such a system is an important step in the modeling of various types of astrophysical phenomena (e. g., Ipser & Lindblom 1991a). The purpose of this paper is to explore some of the consequences arising from the presence of weak viscosity, taking into consideration that its effect is two-fold: a) it changes the value of the unperturbed fluid variables that characterize the equilibrium configuration, and b) it changes the form of the differential equations that govern the evolution of the perturbations. However, we assume that step (a) has already been carried out. We will work within the Cowling approximation (i.e., neglecting gravitational perturbations), which reduces the system to a single differential equation. This approximation is often appropriate (Cowling 1941; Ipser & Lindblom 1992), and is certainly valid for most accretion disks.
In section 2 we present a derivation of the governing partial differential equation. In section 3 we use this result to obtain a formula for the viscosity-induced shift of the complex eigenfrequencies of the perturbations. This enables us, in section 4, to compare our results with those of Ipser & Lindblom (1991a), who perform the calculation for the simpler case of pure uniform rotation of the equilibrium model (which then can contain no effects of viscosity). In section 5, we apply our formalism to thin accretion disks. Explicit integral expressions (involving the eigenfunctions) are obtained for the contributions to the growth (or damping) rate of a mode. We then evaluate the order of magnitude of that rate for the fundamental modes. Except where indicated, we shall employ the notation and conventions of IL.
## 2 GENERALIZED IPSER–LINDBLOM FORMALISM
The equations of motion of a Newtonian fluid are
$`_t\rho +_a(\rho v^a)`$ $`=`$ $`0,`$ (2-1)
$`\rho (_tv^a+v^b_bv^a)`$ $`=`$ $`^ap+\rho ^a\mathrm{\Phi }+_b\sigma ^{ab}.`$ (2-2)
In the cases we shall consider, $`\mathrm{\Phi }`$ represents a specified fixed, axisymmetric gravitational potential (which does not necessarily obey the Newtonian field equation). The viscous stress tensor $`\sigma ^{ab}`$ is defined by
$$\sigma ^{ab}\eta (^av^b+^bv^a)+g^{ab}(\zeta \frac{2}{3}\eta )_cv^c.$$
(2-3)
In this initial investigation, we adopt a generalized barotropic equation of state which specifies the pressure $`p`$ as well as the shear and bulk viscosity coefficients $`\eta `$ and $`\zeta `$ as functions of the mass density $`\rho `$ only. In these equations, $`_t`$ and $`_a`$ represent the partial derivative with respect to time and the spatial Euclidean covariant derivative. Below, $`_a={}_{,a}{}^{}=/x^a`$ indicate spatial partial derivatives, with $`x^a=r,\varphi ,z`$. The Euclidean metric $`g_{ab}`$ and its inverse $`g^{ab}`$ are used to raise and lower tensor indices.
We assume that the viscosity coefficients are small enough that the effects of the viscous term in equation (2-2) can be treated using standard perturbation techniques. To be explicit, we assume that the term $`_b\sigma ^{ab}`$ has a relative magnitude of order $`\alpha 1`$ compared to the dominant terms in the equation. This implies that the viscosity coefficients have a magnitude of order $`\alpha \rho _{}L_{}^2\mathrm{\Omega }_{}`$, where $`\rho _{}`$, $`L_{}`$, and $`\mathrm{\Omega }_{}`$ refer to the typical density, length, and frequency scales of the system. (For an accretion disk, $`L_{}`$ corresponds to the thickness of the disk.) In this way we can treat $`\alpha `$ as a constant perturbation parameter with respect to the inviscid ($`\alpha =0`$) case, and work to first order in $`\alpha .`$
The presence of viscosity will cause a change in the equilibrium values of the dynamical variables. In particular, it can induce meridional components ($`r,z`$) of the velocity field. Using the noncoordinate basis of orthonormal cylindrical coordinates (with unit vectors $`r^a,\varphi ^a`$, and $`z^a`$; and $`g_{ij}=\delta _{ij}`$), the latter can be expressed as
$$v^a=r\mathrm{\Omega }\varphi ^a+v^rr^a+v^zz^a.$$
(2-4)
The last two terms on the right-hand-side of this equation are assumed to have a relative magnitude of order $`\alpha `$ with respect to the first term. (For instance, this is true for a thin accretion disk.)
The equations for the evolution of the Eulerian perturbations (in the Cowling approximation) are
$$_t\delta \rho +v^a_a\delta \rho +\delta \rho _av^a+_a(\rho \delta v^a)=0$$
(2-5)
and
$`_t\delta v^a+v^b_b\delta v^a+\delta v^b_bv^a`$ $`=`$ $`{\displaystyle \frac{^a\delta p}{\rho }}+{\displaystyle \frac{\delta \rho ^ap}{\rho ^2}}{\displaystyle \frac{\delta \rho _b\sigma ^{ab}}{\rho ^2}}`$ (2-6)
$`+\rho ^1\{_b[\eta (^a\delta v^b+^b\delta v^a)]+_b[\delta \eta (^av^b+^bv^a)]`$
$`+^a[(\zeta \frac{2}{3}\eta )_c\delta v^c]+^a[(\delta \zeta \frac{2}{3}\delta \eta )_cv^c]\}.`$
Specifying the functions
$$\mathrm{\Gamma }(r,z)\frac{\rho }{p}\frac{dp}{d\rho },$$
(2-7)
$$\mu (r,z)\frac{\rho }{\eta }\frac{d\eta }{d\rho }$$
(2-8)
effectively closes the system of differential equations. (It turns out that the contributions of $`\delta \zeta `$ are of higher order in $`\alpha `$.) The inviscid case has been studied in an elegant fashion by IL in more generality. (They included gravitational perturbations and allowed a general equation of state, but neglected meridional velocities.) The purpose of this paper is to explore the effects of viscosity on their results under the restrictions of the Cowling approximation and our generalized barotropic equations of state.
We look for normal mode solutions to the equations, which (because of the stationarity and axisymmetry of the equilibrium configuration) have a time dependence of the form $`e^{i\sigma t}`$ and an azimuthal angular dependence of the form $`e^{im\varphi }`$. The constant $`\sigma `$ is the inertial frequency of the mode and the axial mode number $`m`$ is an integer. Let us express equation (2-6) in a way that shows which part depends explicitly on $`\alpha `$ and which does not:
$$(L_b^a+\alpha \mathrm{\Delta }L_b^a)\delta v^b=(R^a+\alpha \mathrm{\Delta }R^a)\delta p.$$
(2-9)
The derivation of the $`\alpha =0`$ terms, $`R^a`$ and $`L_b^a`$, is found in IL. Here we merely quote their results. The operator $`R^a`$ is equal to $`^a\rho ^1`$, where throughout this paper the operator $`^a`$ is to affect all the terms appearing to its right. The operator $`L_b^a`$ is related to the operator $`Q_b^a`$ of IL by $`L_b^a=i(Q^1)_b^a`$, where (in our orthonormal cylindrical coordinates$`r,\varphi ,z`$)
$$Q=\frac{1}{\omega (\omega ^2\kappa ^2)}\left[\begin{array}{ccc}\omega ^2& 2i\mathrm{\Omega }\omega & r(\mathrm{\Omega }^2)_{,z}\\ i\omega \kappa ^2/2\mathrm{\Omega }& \omega ^2& ir\omega \mathrm{\Omega }_{,z}\\ 0& 0& \omega ^2\kappa ^2\end{array}\right].$$
(2-10)
The corotating frequency is $`\omega =\sigma +m\mathrm{\Omega }(r,z)`$, and $`\kappa ^2=2\mathrm{\Omega }(2\mathrm{\Omega }+r\mathrm{\Omega }_{,r})`$ is the square of the radial epicyclic frequency. We note that $`\omega (\omega ^2\kappa ^2)`$ is the determinant of $`Q^1`$. (Also note that IL use a cylindrical coordinate basis and that their roles of $`\sigma `$ and $`\omega `$ are switched with respect to ours.) The derivation of the expressions for $`\alpha \mathrm{\Delta }L_b^a`$ and $`\alpha \mathrm{\Delta }R^a`$ is straightforward and can be found in the Appendix.
The key to the procedure is the fact that $`Q`$ is purely algebraic. This allows us to solve for the velocity perturbation:
$$\delta v^c=(L^1)_a^c(R^a+\alpha \mathrm{\Delta }R^a)\delta p\alpha (L^1)_a^c\mathrm{\Delta }L_b^a\delta v^b.$$
(2-11)
Thus, to the same order in $`\alpha `$,
$$\delta v^c(L^1)_a^c(R^a+\alpha \mathrm{\Delta }R^a)\delta p\alpha (L^1)_a^c\mathrm{\Delta }L_b^a(L^1)_d^bR^d\delta p.$$
(2-12)
Terms of order $`\alpha ^2`$ and higher have been dropped (and will be ignored from now on).
Taking into account equation (2-7), we can express equation (2-5) in the following way:
$$(\omega i\alpha 𝒞)(\rho /p\mathrm{\Gamma })\delta pi_a(\rho \delta v^a)=0,$$
(2-13)
where the operator
$$\alpha 𝒞r^1(rv^r)_{,r}+v_{,z}^z+v^r_r+v^z_z.$$
(2-14)
Throughout this paper, the operators $`_r`$ ($`=/r`$) and $`_z`$ ($`=/z`$) are to affect all the terms appearing to their right. Substituting equation (2-12) in equation (2-13), we obtain
$$_a(\rho Q_b^a^b\rho ^1\delta p)+\frac{\omega \rho }{p\mathrm{\Gamma }}\delta p=\alpha \left[_a(\rho Q_b^a\mathrm{\Delta }R^b\delta p)i_a(\rho Q_b^a\mathrm{\Delta }L_c^bQ_d^c^d\rho ^1\delta p)+i𝒞\frac{\rho }{p\mathrm{\Gamma }}\delta p\right].$$
(2-15)
We have thus reduced the original problem to a single partial differential equation for one unknown, $`\delta p`$. Once solved, $`\delta \rho `$, $`\delta \eta `$, and $`\delta v^a`$ can be obtained via equations (2-7), (2-8), and (2-12), respectively. Our master equation (2-15) plays the same role as equation (22) of IL, which has the same structure within the Cowling approximation.
## 3 EIGENFREQUENCY SHIFTS
Following IL, for reasons that will become apparent below we will use a new variable:
$$\delta V\delta p/\rho \omega .$$
(3-1)
With this definition we can express equation (2-15) as
$$(+\alpha _\mathrm{v})\delta V=0,$$
(3-2)
where
$``$ $``$ $`_a\rho Q_b^a^b\omega +\omega ^2\rho ^2/p\mathrm{\Gamma },`$ (3-3)
$`_\mathrm{v}`$ $``$ $`_a\rho Q_b^a\mathrm{\Delta }R^b\omega \rho +i_a\rho Q_b^a\mathrm{\Delta }L_c^bQ_d^c^d\omega i𝒞\omega \rho ^2/p\mathrm{\Gamma }.`$ (3-4)
The operators $``$ and $`_\mathrm{v}`$ depend upon $`\alpha `$ implicitly through the ($`r`$ and $`z`$ dependent) properties $`v^a,\rho ,p,\eta ,\zeta `$ of the unperturbed model as well as through the eigenfrequency $`\sigma `$. However, we assume that the equilibrium model already incorporates the effects of viscosity, at least through first order in $`\alpha `$. Thus we shall be concerned only with the effects of viscosity on the perturbations, which then changes the operator $``$ only through its dependence on $`\sigma `$.
Expanding the eigenfrequency $`\sigma =\sigma _0+\mathrm{\Delta }\sigma +\mathrm{}`$ and the eigenfunction $`\delta V=\delta V_0+\alpha \delta V_1+\mathrm{}`$, we can express equation (3-2) in the following way, to first order in $`\alpha `$:
$$\left[+(/\sigma )\mathrm{\Delta }\sigma +\alpha _\mathrm{v}\right]_{\sigma _0}(\delta V_0+\alpha \delta V_1)=0.$$
(3-5)
All operators are now evaluated at the adiabatic eigenfrequency $`\sigma _0`$. Using the fact that $`\delta V_0=0`$, we then obtain
$$\alpha \delta V_1+\mathrm{\Delta }\sigma (/\sigma )\delta V_0+\alpha _\mathrm{v}\delta V_0=0.$$
(3-6)
It turns out that with the boundary condition $`p=0`$ on a smooth, compact, finite area surface of the body, $``$ is Hermitean (IL). This implies that
$$\delta V_0^{}\delta V_1d^3x=(\delta V_0)^{}\delta V_1d^3x=0.$$
(3-7)
Thus if we multiply equation (3-6) by $`\delta V_0^{}`$ and integrate, we obtain
$$\mathrm{\Delta }\sigma /\sigma +\alpha _\mathrm{v}=0;$$
(3-8)
where
$$𝒪\delta V_0^{}𝒪\delta V_0d^3x,$$
(3-9)
only for operators $`𝒪`$ evaluated at $`\sigma =\sigma _0`$.
Then we find (from the structure of $``$ and $`_\mathrm{v}`$) that
$$\mathrm{Re}(\mathrm{\Delta }\sigma )=0,i\mathrm{Im}(\mathrm{\Delta }\sigma )=\frac{\alpha _\mathrm{v}}{/\sigma }.$$
(3-10)
Equation (3-10) allows us to compute the viscosity-induced shift of the eigenfrequencies, provided we are given an equilibrium model. Equation (3-6) can then be used to obtain the eigenfunction correction $`\delta V_1`$.
## 4 PERTURBATIONS OF UNIFORMLY ROTATING BODIES
In order to compare our results with those obtained by Ipser & Lindblom (1991a), we need to invoke their assumptions of uniform rotation and vanishing of meridional velocity (i.e., we need to set $`\mathrm{\Omega }_{,r}=\mathrm{\Omega }_{,z}=v^r=v^z=0`$). Thus, there are no effects of viscosity on their equilibrium model since the viscous shear tensor vanishes. In addition, since we do not consider self-gravitational effects, we need to set $`\delta \mathrm{\Phi }=0`$ in their formalism. Under these assumptions, equation (3-10) reproduces their results \[equations (27), (28), and (29)\] for the imaginary part of the frequency induced by viscous effects. \[Their $`\tau ^1`$ corresponds to our Im($`\mathrm{\Delta }\sigma `$), and their canonical energy $`E`$, to our quantity $`(\omega /2)/\sigma `$.\] When dealing with differentially rotating bodies, our approach is more straightforward than an extension of theirs, which employs a Lagrangian approach.
## 5 APPLICATION TO ACCRETION DISKS
The results obtained so far can be applied to the special case of a thin accretion disk. Its physical properties allow us to retain relatively few terms from the ones contained in equation (3-10). Adiabatic oscillations of such accretion disks have been studied within a radial WKB approximation (Nowak & Wagoner 1992; Perez et al. 1997; Silbergleit & Wagoner 1999) $`\lambda _rr_{}`$, with $`\lambda _rh_{}`$ for most of the low-lying modes. The length scales $`2h_{}`$ and $`r_{}`$ are the typical thickness and radial size of the disk region of interest, and $`\lambda _r`$ is the radial scale of variation of the eigenfunction. In addition, we shall only consider values of the axial mode number $`|m|1`$.
The system is characterized by the small parameter $`ϵh_{}/r_{}`$. In order to use our equations we will also need a model for the unperturbed disk, which we take to be similar to the one developed by Kita & Kluźniak (1998). In agreement with their results, we take $`v^r\alpha ϵ^2r\mathrm{\Omega }`$ and $`v^zϵv^r`$.
Taking into account these considerations, equation (3-10) gives (to lowest order in $`\alpha `$ and $`ϵ`$) the viscous damping rate
$$\frac{1}{\tau }\text{Im}(\mathrm{\Delta }\sigma )=\frac{(F_1+F_2+F_3+F_4)d^3x}{F_0d^3x}.$$
(5-1)
To obtain this expression, we employed (many) integrations by parts, with the boundary conditions that $`\delta V_0`$, $`\eta `$, and $`\zeta `$ vanish. The integrands are
$`F_0`$ $`=`$ $`2\omega \rho \left({\displaystyle \frac{\rho }{\mathrm{\Gamma }p}}|\delta V_0|^2+\kappa ^2S^2\left|{\displaystyle \frac{\delta V_0}{r}}\right|^2\right),`$ (5-2)
$`F_1`$ $`=`$ $`\omega [\omega ^2(2\eta +\xi )+\kappa ^2\eta ]\left|{\displaystyle \frac{}{r}}\left(S{\displaystyle \frac{\delta V_0}{r}}\right)\right|^2,`$ (5-3)
$`F_2`$ $`=`$ $`\omega ^1(2\eta +\xi )\left|{\displaystyle \frac{^2\delta V_0}{z^2}}\right|^2,`$ (5-4)
$`F_3`$ $`=`$ $`\omega [\eta (\omega ^2+S^2(\omega ^2+\kappa ^2))+2S(\eta +\xi )]\left|{\displaystyle \frac{^2\delta V_0}{rz}}\right|^2,`$ (5-5)
$`F_4`$ $`=`$ $`\omega S\left[(\mu 1){\displaystyle \frac{r\eta \rho }{\mathrm{\Gamma }p}}\left({\displaystyle \frac{\mathrm{\Omega }^2}{r}}\right)+{\displaystyle \frac{^2\xi }{z^2}}\right]\left|{\displaystyle \frac{\delta V_0}{r}}\right|^2.`$ (5-6)
We have introduced $`S1/(\omega ^2\kappa ^2)`$ and $`\xi \zeta \frac{2}{3}\eta `$. The IL operator defined by equation (3-3), which reduces to the form
$$=\frac{\omega ^2\rho ^2}{\mathrm{\Gamma }p}+\omega ^2\rho \frac{}{r}\left(S\frac{}{r}\right)+\frac{}{z}\left(\rho \frac{}{z}\right),$$
(5-7)
has also been employed. It is interesting to note that the terms proportional to $`v^r`$ and $`v^z`$ do not appear in this limit, since their contribution is of higher order in $`ϵ`$. This result justifies their neglect in previous analyses of thin disks.
Let us now consider in more detail the effect of viscosity on observationally relevant fundamental modes of relativistic accretion disks (near a black hole or neutron star). Since we are employing a Newtonian analysis, but use the properties of the relativistic eigenfunctions $`\delta V_0`$, our results are only approximate. The trapping of the g- and p-modes is produced by the fact that relativistic effects make the radial epicyclic frequency $`\kappa `$ less than the angular velocity of free-particle circular orbits. This reduction produces an inner edge of the disk where $`\kappa (r)`$ vanishes, with the maximum value of $`\kappa (r)`$ at a slightly larger radius. The trapping of the c-modes is related to the reduction of the vertical epicyclic frequency below the angular velocity, produced by the angular momentum of the central mass.
* The g-modes are trapped where the comoving frequency $`|\omega |<\kappa `$, near the maximum of $`\kappa (r)`$. They are mathematically similar to the internal gravity (buoyancy) modes of stars. For the fundamental ($`m=0`$) g-mode (Perez et al. 1997), $`_rϵ^{1/2}_z`$, $`_z1/h_{}`$, and $`_\varphi =0`$ when applied to $`\delta V_0`$. Single terms in $`F_0`$ and $`F_3`$ dominate, giving
$$\frac{1}{\tau }\frac{\eta (\sigma ^2+\kappa ^2)S^2|_r_z\delta V_0|^2d^3x}{2\rho \kappa ^2S^2|_r\delta V_0|^2d^3x}\alpha \mathrm{\Omega }_{},$$
(5-8)
with $`S1/(ϵ\sigma ^2)`$. The sign of this result agrees with that obtained by Nowak & Wagoner (1992), although the magnitude of their growth rate was greater.
* The (inner) p-modes are trapped between the inner edge of the disk and the increasing portion of $`\kappa (r)`$, with $`|\omega |>\kappa `$. They are similar to the pressure (acoustic) modes of stars. For the fundamental ($`m=0`$) p-mode (Nowak & Wagoner 1991), $`_rϵ^{1/3}_z`$, $`_z1/h_{}`$, and $`_\varphi =0`$ when applied to $`\delta V_0`$. The numerator is dominated by the same term in $`F_3`$ as above, as well as $`F_2`$, giving
$$\frac{1}{\tau }\frac{[(2\eta +\xi )\sigma ^2|_z^2\delta V_0|^2+\eta (\sigma ^2+\kappa ^2)S^2|_r_z\delta V_0|^2]d^3x}{2[\rho c_s^2|\delta V_0|^2+\rho \kappa ^2S^2|_r\delta V_0|^2]d^3x}\frac{\alpha \mathrm{\Omega }_{}}{(\sigma /\mathrm{\Omega }_{})^2+ϵ^{2/3}},$$
(5-9)
with $`c_s^2=\mathrm{\Gamma }p/\rho `$ and $`S1/\sigma ^2`$. The first (second) integral in both the numerator and the denominator dominates when $`ϵ`$ is smaller (greater) than $`(\sigma /\mathrm{\Omega }_{})^31`$. The sign of this result agrees with that obtained in a local analysis by Kato et al. (1998), although its form differs somewhat from theirs.
* The c-modes are trapped between the inner edge of the disk and the radius where the Lense–Thirring frequency is the eigenfrequency. For the fundamental ($`m=\pm 1`$) c-mode (Silbergleit & Wagoner 1999), $`_\varphi =\pm 1`$, $`_z1/h_{}`$ but $`_z^2=0`$ when applied to $`\delta V_0`$. Thus this mode is vertically incompressible. One finds that
$$\frac{1}{\tau }\frac{(F_3+F_4)d^3x}{2\rho c_s^2\omega |\delta V_0|^2d^3x}\pm \alpha \left(\frac{h_{}}{\lambda _r}\right)^2\mathrm{\Omega }_{},$$
(5-10)
where $`S\omega ^2`$. Note that this mode may either grow or be damped by viscosity.
We remind the reader that the real part of the corrections to all eigenfrequencies vanish. Future generalizations of this work should include the effects of buoyancy (nonbarotropic equation of state) and effects of higher order in $`ϵh_{}/r_{}`$.
This research was supported in part by NASA grants NAG 5-3102 to R.V.W. and NAS 8-39225 to Gravity Probe B. We thank John Friedman, Lee Lindblom, and Alex Silbergleit for helpful suggestions.
## Appendix A APPENDIX: Determination of $`\mathrm{\Delta }R^a`$ and $`\mathrm{\Delta }L_b^a`$
The components of $`_b\sigma ^{ab}`$ in orthonormal cylindrical coordinates are
$`_b\sigma ^{rb}`$ $`=`$ $`2(\eta v_{,r}^r)_{,r}+r^1[\eta (r^1v_{,\varphi }^r+v_{,r}^\varphi r^1v^\varphi )]_{,\varphi }+[\eta (v_{,z}^r+v_{,r}^z)]_{,z}`$ (A1a)
$`+2\eta r^1(v_{,r}^rr^1v^rr^1v_{,\varphi }^\varphi )+\{(\zeta \frac{2}{3}\eta )[r^1(rv^r)_{,r}+r^1v_{,\varphi }^\varphi +v_{,z}^z]\}_{,r};`$
$`_b\sigma ^{\varphi b}`$ $`=`$ $`[\eta (v_{,r}^\varphi r^1v^\varphi +r^1v_{,\varphi }^r)]_{,r}+2\eta r^1(v_{,r}^\varphi r^1v^\varphi +r^1v_{,\varphi }^r)+2r^2[\eta (v^r+v_{,\varphi }^\varphi )]_{,\varphi }`$ (A1b)
$`+[\eta (v_{,z}^\varphi +r^1v_{,\varphi }^z)]_{,z}+r^1\{(\zeta \frac{2}{3}\eta )[r^1(rv^r)_{,r}+r^1v_{,\varphi }^\varphi +v_{,z}^z]\}_{,\varphi };`$
$`_b\sigma ^{zb}`$ $`=`$ $`[\eta (v_{,z}^r+v_{,r}^z)]_{,r}+\eta r^1(v_{,z}^r+v_{,r}^z)+2(\eta v_{,z}^z)_{,z}+\eta r^1(v_{,z}^\varphi +r^1v_{,\varphi }^z)]_{,\varphi }`$ (A1c)
$`+\{(\zeta \frac{2}{3}\eta )[r^1(rv^r)_{,r}+r^1v_{,\varphi }^\varphi +v_{,z}^z]\}_{,z}.`$
These expressions and their perturbations allow us to compute the operator $`_\mathrm{v}`$. There are two terms in the right-hand-side of equation (2-6) that contribute to $`\alpha \mathrm{\Delta }R`$. We will consider each in order. The first contribution comes from $`\rho ^2\delta \rho _b\sigma ^{ab}`$. The order $`\alpha `$ part of this contribution can be obtained from the above equations, noticing that $`_\varphi =0`$ when acting on an equilibrium quantity and that the variables $`v^r,v^z,\eta ,\zeta ,\delta \eta `$, and $`\delta \zeta `$ are all of order $`\alpha `$:
$$\alpha \mathrm{\Delta }R_A=\frac{1}{\rho ^2}\left[\begin{array}{c}0\\ (r\rho /p\mathrm{\Gamma })(\eta \mathrm{\Omega }_{,z})_{,z}+(_rr+2)(\eta \rho \mathrm{\Omega }_{,r}/p\mathrm{\Gamma })\\ 0\end{array}\right].$$
(A2)
The second contribution comes from $`\rho ^1_b[\delta \eta (^av^b+^bv^a)]`$:
$$\alpha \mathrm{\Delta }R_B=\frac{1}{\rho }\left[\begin{array}{c}im\eta \mu \mathrm{\Omega }_{,r}/p\mathrm{\Gamma }\\ r_z(\eta \mu \mathrm{\Omega }_{,z}/p\mathrm{\Gamma })+(_rr+2)(\eta \mu \mathrm{\Omega }_{,r}/p\mathrm{\Gamma })\\ im\eta \mu \mathrm{\Omega }_{,z}/p\mathrm{\Gamma }\end{array}\right].$$
(A3)
The term $`\rho ^1^a[(\delta \zeta \frac{2}{3}\delta \eta )_cv^c]`$ is of higher order in $`\alpha `$ and can therefore be ignored. Thus
$$\alpha \mathrm{\Delta }R=\alpha \mathrm{\Delta }R_A+\alpha \mathrm{\Delta }R_B.$$
(A4)
For the determination of $`\alpha \mathrm{\Delta }L`$, let us start with the contribution from the left-hand-side of equation (2-6). The components of the latter are
$`_t\delta v^r+v^b_b\delta v^r+\delta v^b_bv^r`$ $`=`$ $`i\sigma \delta v^r+im\mathrm{\Omega }\delta v^r2\mathrm{\Omega }\delta v^\varphi +v^r\delta v_{,r}^r+v^z\delta v_{,z}^r`$ (A5a)
$`+v_{,r}^r\delta v^r+v_{,z}^r\delta v^z,`$
$`_t\delta v^\varphi +v^b_b\delta v^\varphi +\delta v^b_bv^\varphi `$ $`=`$ $`i\sigma \delta v^\varphi +im\mathrm{\Omega }\delta v^\varphi +\mathrm{\Omega }\delta v^r+v_{,r}^\varphi \delta v^r+v_{,z}^\varphi \delta v^z`$ (A5b)
$`+v^r\delta v_{,r}^\varphi +v^z\delta v_{,z}^\varphi +r^1v^r\delta v^\varphi ,`$
$`_t\delta v^z+v^b_b\delta v^z+\delta v^b_bv^z`$ $`=`$ $`i\sigma \delta v^z+im\mathrm{\Omega }\delta v^z+v^r\delta v_{,r}^z+v^z\delta v_{,z}^z+v_{,r}^z\delta v^r`$ (A5c)
$`+v_{,z}^z\delta v^z.`$
The first contribution to $`\alpha \mathrm{\Delta }L`$ comes from the terms with $`v^r`$ or $`v^z`$:
$$\alpha \mathrm{\Delta }L_A=(v^r_r+v^z_z)𝐈+\left[\begin{array}{ccc}v_{,r}^r& 0& v_{,z}^r\\ 0& v^r/r& 0\\ v_{,r}^z& 0& v_{,z}^z\end{array}\right],$$
(A6)
where I is the identity matrix. Let us turn now to the contributions to $`\alpha \mathrm{\Delta }L`$ coming from the right-hand-side of equation (2-6). The first of such contributions comes from the term $`\rho ^1_b[\eta (^a\delta v^b+^b\delta v^a)]`$:
$$\alpha \mathrm{\Delta }L_B=\frac{1}{\rho }[\begin{array}{cc}2_r\eta _r+2\eta _rr^1m^2r^2\eta +_z\eta _z& im\eta r^1(_r3r^1)\\ (_r+4r^1)im\eta r^1& (_rr+2)\eta _rr^12m^2r^2\eta +_z\eta _z\\ (_r+r^1)\eta _z& im\eta r^1_z\end{array}$$
$$\begin{array}{c}_z\eta _r\\ imr^1_z\eta \\ (_r+r^1)\eta _r+2_z\eta _zm^2r^2\eta \end{array}].$$
(A7)
The other contribution comes from the term $`\rho ^1^a[(\zeta \frac{2}{3}\eta )_c\delta v^c]`$:
$$\alpha \mathrm{\Delta }L_C=\frac{1}{\rho }\left[\begin{array}{ccc}_r\xi r^1_rr& im_r\xi r^1& _r\xi _z\\ im\xi r^2_rr& m^2r^2\xi & im\xi r^1_z\\ _z\xi r^1_rr& imr^1_z\xi & _z\xi _z\end{array}\right],$$
(A8)
where $`\xi \zeta \frac{2}{3}\eta `$. Thus
$$\alpha \mathrm{\Delta }L=\alpha \mathrm{\Delta }L_A+\alpha \mathrm{\Delta }L_B+\alpha \mathrm{\Delta }L_C.$$
(A9) |
warning/0002/hep-th0002247.html | ar5iv | text | # A UNIFORM MODEL OF THE MASSIVE SPINNING PARTICLE IN ANY DIMENSION.
## 1 Introduction
The Lagrangian description of the relativistic spinning particles is one of the recurrently discussed themes in high energy physics, having a long history. The retrospective exposition of the question and some basic references can be found in the review . In the context of $`M`$-theory the models of spinning particles awake today fresh interest as special but highly nontrivial examples of $`0`$-branes which along with some other extended objects are considered to be the basic ingredients of the non-perturbative string theory . Since the target space of a consistent string theory has certainly more than four dimensions ($`10`$ or perhaps $`11`$) the observable space-time is supposed to result from the Kaluza-Klein compactification of extra dimensions in the low energy limit. In so doing the effective dimension of p-brane may decrease down to zero when it is considered from the viewpoint of four-dimensional observer. The classical example is the double dimensional reduction of $`d=11`$ supermembrane to the type II A string in $`d=10`$ . This opens up an interesting possibility to interpret the spinning particles as low energy effective models of the p-branes in higher-dimensional space-time, which compact directions are associated with the spinning degrees of freedom of the particle.
All this gives rise to the question about the construction of the mechanical models of relativistic spinning particles in the space-time of arbitrary dimension. The reach kinematical symmetries underlying the models enable one to treat them as elementary dynamical systems in the Souriau sense and to apply for their description the full machinery of the symplectic geometry. (For applications of this approach see also ). In the framework of this scheme the whole dynamical information about the space-time and phase-space evolution of the system is encoded in a presymplectic manifold $``$ being a homogeneous transformation space of the group $`G`$ for which the system is elementary one. For a given physical system, the choice of $``$ is very ambiguous and there is no precise prescription for it. Fortunately, there is no matter how to choose particular $``$ when describing free particle: any $``$ leads to the proper classical dynamics<sup>a</sup><sup>a</sup>aFor example, one can always identify $``$ with the underlying symmetry group $`G`$ itself..
Various models of spinning particles are known with $`4d`$ symmetry groups: Poincaré , de Sitter and Galilei . The $`3d`$ and $`6d`$ analogues to these models can be found in ref. , for superextensions see . The covariant operatorial quantization of these models leads to the Hilbert space of physical states carrying the unitary irreducible representation of the respective groups. Some of these models have, however, a common problem of constructing a consistent extension to the case of spinning particle subject to exterior fields. The obstacle is that the space-time dynamics of the particle, arising in these models in some cases, is characterized by two-dimensional world-tubes rather then world-lines. The non-local behavior of such a type, sometimes referred to as the phenomenon of Zitterbewegung , is usual for the relativistic particles with spin and presents the main obstruction to the switching on a local interaction. Note that the difficulty is not inherent to these systems, as they are, but it is rather related to the way of their description, which basically depends on the choice of $``$. In particular, the model of the ref , in $`d=4`$, does not display the Zitterbewegung, and it allows the consistent interaction.
In recent paper we have constructed the model for a massive particle of integer spin living in $`d`$–dimensional space-time and coupled to an arbitrary background of gravity and electromagnetism. The underlying presymplectic manifold was identified with that for the spinless particle times a regular (co)adjoint orbit of space rotations group
$$\begin{array}{c}=_{spinless}\times 𝒪_𝐬,\\ \\ _{spinless}=𝐑^{d1,1}\times B,𝒪_𝐬=SO(d1)/[SO(2)]^r\end{array}$$
(1)
Here $`r=rankSO(d1)=[(d1)/2]`$ and $`B`$ stands for the upper sheet of the mass hyperboloid: $`p_Ap^A=m^2,p_0>0`$. The regularity of $`𝒪_𝐬`$ means that the whole space of invariant presymplectic structures on $``$ is parametrized by $`r+1`$ numbers $`m,𝐬=(s_1,s_2,\mathrm{}s_r)`$ associated with mass and spin(s) of the particle <sup>b</sup><sup>b</sup>bThis is just the number of parameters labelling a general massive representation of Poincare’ group. Indeed, fixing mass of the particle reduces the classification problem for the Poincare’ group representations to the one for the Wigner little group $`SO(D1)`$. According to the general Borel-Bott-Weil theory , there is a one-to-one correspondence between regular (co-)adjoint orbits of the orthogonal groups and theirs representations. The use of the other (irregular) orbits would lead only to the special spin representations, which can be obtained in our model by specifying values of $`𝐬^{}`$.. The space-time motion of the particle is described here by the one-dimensional world-lines (time-like geodesics of Minkowski space) and, as a result, the model is free from the above mentioned obstacle to the interaction. For the sake of explicit Poincaré covariance, the orbit $`𝒪_𝐬`$ was symplectically embedded into $`_{i=1}^r𝐑_𝐂^{d1,1}`$ equipped with the natural action of the Lorentz group $`SO(d1,1)`$ and invariant symplectic form. This realization for $`𝒪_𝐬`$ proves to be especially suitable for the covariant quantization of integer-spin particle but it becomes inadequate when trying to consider half-integer spins since the quantum-mechanical description for the latter case is based on the group $`Spin(d1,1)`$ rather than $`SO(d1,1)`$. In order to take into account the half-integer spins, one should replace, from the very beginning, the proper Lorentz group $`SO(d1,1)`$ by its double covering $`Spin(d1,1)`$.
In this paper we propose the new construction for the spinning sector of the massive particle which provides a uniform quantum-mechanical description for both integer and half-integer spins. For these ends, the phase space of spinning degrees of freedom $`𝒪_𝐬`$ is embedded into the carrier space of the Lorentz group $`𝐂^{2^{[d/2]}}_{i=1}^{r1}𝐑_𝐂^{d1,1}`$, where the first factor transforms under the spinor representation. Then the (half-)integer spin representations may be obtained by applying either geometric or Dirac quantization to this model. A remarkable property of this realization for $`𝒪_𝐬`$ is that the use of spinor variables makes possible to resolve the mass-shell condition $`p_Ap^A=m^2`$ in a Lorentz-invariant manner. It can be thought about as an extension to the massive case of the twistor realization known for the isotropic momenta of a massless particle in special dimensions: 3, 4, 6 and 10. The distinction is that the spinor variable carries now an information about both the mass hyperboloid $`B`$ and the internal space for spin $`𝒪_𝐬`$. The details of this construction are presented in the next section together with the generalization to the case of minimal coupling to exterior gravitational and electromagnetic fields.
In Section 3, we reformulate the model as Hamiltonian system with the first and second class constraints and study the physical spectrum of the theory within the Dirac quantization scheme. The physical wave functions, being extracted by the quantum operators assigned to the constraints, are shown to be in one-to-one correspondence with the Poincaré-irreducible (spin-)tensor fields in Minkowski space.
We conclude the paper by discussing the obtained results and some further perspectives.
## 2 Classical description
As it was mentioned in the previous section the extended phase space of the massive spinning particle is constructed to be the direct product of the extended phase space of the spinless particle and the manifold $`𝒪_𝐬`$, responsible for spinning degrees of freedom. The former factor is standardly embedded into the cotangent bundle $`T^{}(𝐑^{d1,1})`$ of the Minkowski space by the mass-shell condition
$$p_Ap^A+m^2=0$$
(2)
while the latter is identified with so-called flag manifold, which points may be viewed as the sequences
$`0`$ $``$ $`V_1V_2\mathrm{}V_r𝐑_𝐂^{d1,1},`$
$`p`$ $``$ $`V_k,dimV_k=k`$
of complex $`p`$-transversal vector subspaces of complexified Minkowski space, telescopically embedded into each other. In the previous paper , we have suggested the holomorphic parametrization for $`𝒪_𝐬`$ with the help of $`r`$ independent complex vectors $`Z_i^A𝐑_𝐂^{d1,1},i=1,2,\mathrm{}r;A=0,1,\mathrm{}d1`$ subject to the conditions
$$(Z_i,Z_j)=0,(p,Z_i)=0,Z_i^AZ_j^A\mathrm{\Lambda }_i^j,$$
(3)
where $`\mathrm{\Lambda }_i^j`$ is a complex non-degenerate upper-triangular $`r\times r`$ matrix, and $`(\mathrm{},\mathrm{})`$ denotes the inner product with respect to the Minkowski metric $`\eta _{AB}=diag(,+\mathrm{},+)`$. In this realization, each $`V_k`$ is spanned by the vectors $`Z_1^A,Z_2^A,\mathrm{}Z_k^A`$ and the last relation in (3) establishes the equivalence between all such frames in $`V_k`$.
To construct a spinor realization for $`𝒪_𝐬`$, it is sufficient to parametrize a subspace $`V_r`$ by a spinor variable. Consider the commuting Dirac spinor $`\psi _a,a=1,\mathrm{}2^{[d/2]}`$ subject to the following constraints and equivalence relations<sup>c</sup><sup>c</sup>cIn certain dimensions we could parametrize the manifold $`𝒪_𝐬`$ by Weyl or Majorana spinors as well. This, however, leads to different constraints, which depend on each the specific dimension. The parametrization by Dirac spinors does not depend on the dimension explicitly, that seems to be more convenient for the uniform description of the spin in higher dimensions.
$$\begin{array}{c}\psi _a\lambda \psi _a,\lambda 𝐂\{0\},\\ \\ p_A\mathrm{\Gamma }_a^{Ab}\psi _b=m\psi _a,Z_{kA}\mathrm{\Gamma }_a^{Ab}\psi _b=0,k=1,\mathrm{}r\end{array}$$
(4)
Since some subsequent expressions may differ for the cases of even- and odd dimensions, the formulae will be labeled with the letters $`a`$ and $`b`$ for the former and latter case respectively. Accounting (4) and the Fierz identities, the spinor bilinear is decomposed in the basis of the Clifford algebra generated by $`\mathrm{\Gamma }`$ -matrices as follows:
$$\psi \stackrel{~}{\psi }=M_{A(r)}\{\mathrm{\Gamma }^{A(r)}\frac{(1)^r}{m}p_B\mathrm{\Gamma }^{BA(r)}\},$$
$`(5.a)`$
$$\psi \stackrel{~}{\psi }=M_{A(r)}\mathrm{\Gamma }^{A(r)},$$
$`(5.b)`$
where $`\stackrel{~}{\psi }`$ is the charge conjugated spinor<sup>d</sup><sup>d</sup>d The charge conjugation matrix $`C`$ is determined from the relation $`\mathrm{\Gamma }_A^T=(1)^rC\mathrm{\Gamma }_AC^1`$ $`\stackrel{~}{\psi }{}_{}{}^{a}=(\psi C)^a`$. Hereafter we use the shorthand notation $`A(r)=A_1\mathrm{}A_r`$. Tensor $`M_{A(r)}=\frac{(1)^{r(r+1)/2}}{2^{[d/2]}r!}(\stackrel{~}{\psi }\mathrm{\Gamma }_{A(r)}\psi )`$ obeys equations
$$\begin{array}{c}M_{[A(r)}Z_{kB]}=0,Z_k^AM_{AA(r1)}=0,k=1,\mathrm{}r,\\ \\ p^AM_{AA(r1)}=0,M_{A(r)}\lambda ^2M_{A(r)}\end{array}$$
(6)
Making use of (6) one can express $`M_{A(r)}`$ in terms of the $`Z_i`$ . For these ends it is convenient to introduce another parametrization of $`V_k`$ with the help of $`k`$-forms $`𝒵^k`$ determined from the equations
$$𝒵^k(\overline{Z}_1,\overline{Z}_2,\mathrm{}\overline{Z}_k)=det(Z_i\overline{Z}_j),i,j=1,\mathrm{}k,$$
(7)
that establish a relation between $`𝒵`$ and $`Z`$
$$𝒵_{A(k)}^kZ_{1[A_1}Z_{2A_2}\mathrm{}Z_{kA_k]}$$
(8)
Tensors $`𝒵_{A(k)}^k`$ satisfy the following relations
$$\begin{array}{c}𝒵_{[A(i)}^i𝒵_{B]B(k1)}^k=0,\eta ^{AB}𝒵_{AA(i1)}^i𝒵_{BB(k1)}^k=0,p^A𝒵_{AA(k1)}^k=0,\\ \\ 𝒵_{A(k)}^ka_k𝒵_{A(k)}^k,i,k=1,\mathrm{}r,ik,a_k𝐂\{0\},\end{array}$$
(9)
resulting from the definitions (3). Now it is easy to see that the general solution to the equations (6) has the form
$$M_{A(r)}𝒵_{A(r)}^r$$
(10)
Thus the $`r`$-dimensional complex vector subspace $`V_r`$ can be parametrized by the spinor $`\psi `$ subject to the conditions (4). Making use of ( 8) we obtain the equivalent parametrization of $`𝒪_𝐬`$ in terms of $`(r1)`$ complex vectors $`Z_i^A,i=1,2,\mathrm{}r1,A=0,1,\mathrm{}d1`$ and the spinor $`\psi `$ subject to equivalence relations
$$Z_i^AZ_j^A\mathrm{\Lambda }_i^j,\psi \lambda \psi ,\lambda 𝐂\{0\}$$
(11)
and constraints
$$(p,Z_i)=0,p_A\mathrm{\Gamma }^A\psi =m\psi ,$$
(12)
$$(Z_i,Z_j)=0,Z_{iA}\mathrm{\Gamma }^A\psi =0$$
(13)
Here $`\mathrm{\Lambda }_i^j`$ is a complex non-degenerate upper-triangular $`(r1)\times (r1)`$ matrix. In terms of introduced objects the most general expression for the Kähler potential on $`𝒪_𝐬`$ is
$$\begin{array}{c}\mathrm{\Phi }=\frac{1}{2}\mathrm{ln}(\mathrm{\Delta }_1^{s_1}\mathrm{}\mathrm{\Delta }_{r1}^{s_{r1}}\mathrm{\Delta }_r^{s_r}),\\ \\ \mathrm{\Delta }_i=Z_{1A_1}\mathrm{}Z_{iA_i}\overline{Z}_1^{[A_1}\mathrm{}\overline{Z}_i^{A_i]},i=1,\mathrm{},r1,\mathrm{\Delta }_r=(\overline{\psi }\psi )^2,\end{array}$$
(14)
where $`\overline{\psi }^a=(\psi ^{}\mathrm{\Gamma }_0)^a`$ is the Dirac conjugated spinor . Note that $`\mathrm{\Phi }`$ depends on $`p`$ implicitly in view of constraints (12). Under transformations (11) $`\mathrm{\Phi }`$ changes to an additive constant
$$\delta _{\mathrm{\Lambda },\lambda }\mathrm{\Phi }=\underset{k=1}{\overset{r1}{}}\mathrm{ln}\left|\mathrm{\Lambda }_1^1\mathrm{\Lambda }_2^2\mathrm{}\mathrm{\Lambda }_k^k\right|^{s_k}+\mathrm{ln}\left|\lambda \right|^{2s_r}$$
(15)
The direct product structure of $``$ allows to introduce the one-form $`\theta `$ being a sum of conventional one-form $`p_Adx^A`$ on $`_{spinless}`$ describing the space-time dynamics of the particle and a one-form on $`𝒪_𝐬`$ governing the spinning dynamics. We will put
$$\theta =p_Adx^A+d\mathrm{\Phi },$$
(16)
where the action of the star operator on the complex one-forms is defined as $`(\alpha _Idz^I+\beta _Id\overline{z}^I)=i(\alpha _Idz^I\beta _Id\overline{z}^I)`$. Notice that $`\theta `$ is invariant under transformations (11) modulo closed one-form and thus the Hamiltonian action for the system may be chosen as
$$S=\underset{\gamma }{}\theta $$
(17)
The extremals of the action (17) coincide with the leaves of $`\mathrm{ker}d\theta `$. By construction $`\mathrm{ker}d\theta `$ is generated by the only vector field $`𝐕=p_A/x_A`$ and hence the tangent vector to the trajectory is proportional to $`𝐕`$. This means
$$\begin{array}{c}\dot{x}^A=\mu p^A,\dot{p}_A=0,\\ \\ \dot{Z}_i^A=0,\dot{\psi }_a=0,(modulotransformations(\text{11}))\end{array}$$
(18)
where $`\mu =\mu (\tau )`$ is an arbitrary function of proper time which is fixed after particular choice of a world-line parametrization. Thus the particle moves along the time-like geodesics in the Minkowski space while the internal degrees of freedom do not evolve.
The interesting feature of the action (17) is that the Dirac equation on $`\psi `$ subject to the rest constraints (11, 12, 13) may be covariantly resolved with respect to $`p_A`$. To solve this equation, we multiply it by $`\psi _c`$ and make use of the decomposition (5) together with Fierz identities for spinor bilinears. This results in relations
$$p^A(\stackrel{~}{\psi }\mathrm{\Gamma }_{AA(r1)}\psi )=0,(\stackrel{~}{\psi }\mathrm{\Gamma }_{AA(r)}\psi )=\frac{(1)^{r+1}(r+1)}{m}p_{[A}(\stackrel{~}{\psi }\mathrm{\Gamma }_{A(r)}\psi ),$$
$`(19.a)`$
$$p^A(\stackrel{~}{\psi }\mathrm{\Gamma }_{AA(r1)}\psi )=0,ϵ_{A(r)BC(r)}p^B(\stackrel{~}{\psi }\mathrm{\Gamma }^{C(r)}\psi )=(1)^{1+r(r+1)/2}r!mi^r(\stackrel{~}{\psi }\mathrm{\Gamma }_{A(r)}\psi )$$
$`(19.b)`$
Contracting the second equation in (19.a) with $`(\stackrel{~}{\psi }\mathrm{\Gamma }^{A(r)}\psi )^{}`$ one can express $`p_A`$ via the spinor variables
$$p_A=\overline{p}_A(\psi ,\psi ^{})(1)^{r+1}m\frac{[(\stackrel{~}{\psi }\mathrm{\Gamma }_{AB(r)}\psi )(\stackrel{~}{\psi }\mathrm{\Gamma }^{B(r)}\psi )^{}+c.c.]}{2\left|(\stackrel{~}{\psi }\mathrm{\Gamma }_{C(r)}\psi )\right|^2}$$
$`(20.a)`$
For the odd-dimensional case the similar trick (with $`(\stackrel{~}{\psi }\mathrm{\Gamma }^{B(r)}\psi )^{}`$ replaced by
$`ϵ^{A(r)CD(r)}(\stackrel{~}{\psi }\mathrm{\Gamma }_{D(r)}\psi )^{}`$ ) yields
$$p_A=\overline{p}_A(\psi ,\psi ^{})m\frac{(1)^{r(r1)/2}i^rϵ_{AB(r)C(r)}(\stackrel{~}{\psi }\mathrm{\Gamma }^{B(r)}\psi )(\stackrel{~}{\psi }\mathrm{\Gamma }^{C(r)}\psi )^{}}{r!\left|(\stackrel{~}{\psi }\mathrm{\Gamma }_{D(r)}\psi )\right|^2}$$
$`(20.b)`$
Since equations (19) form the overdetermined system, the expressions (20) being inserted back into (19) will lead to some consistency conditions which should be imposed on $`\psi `$ besides the holonomic constraints (13). Note that the mass-shell condition (2) reduces to the purely algebraic relation on $`\psi `$ which is valid due to these constraints. Thus we obtain the parametrization of the massive hyperboloid in terms of the spinor $`\psi _a`$ subject to the conditions (11, 13) and (19) (with $`p`$ replaced by $`\overline{p}`$). This resembles to the twistor parametrization of the light-cone <sup>d</sup><sup>d</sup>dFor the models of massless spinning (super)particles which exploit the twistor parametrizations see, e.g. which, however, is connected with the division algebras and therefore exists in $`3`$, $`4`$, $`6`$ and $`10`$ dimensions only. The derived construction may be considered as its generalization to the case of arbitrary dimensional space-time with the exception that the spinor $`\psi `$ contains information about both the space-time momentum and intrinsic degrees of freedom.
Substituting $`\overline{p}_A`$ in (16), we immediately derive the Lagrangian of the system
$$=\overline{p}_A(\psi ,\psi ^{})\dot{x}^Ai(\dot{Z}_i^A_A^i\mathrm{\Phi }+\dot{\psi }_a^a\mathrm{\Phi }c.c.)$$
(21)
The Lagrangian is obviously invariant under the global Poincaré transformations, reparametrizations of the world-line and changes to a total derivative under the gauge transformations associated with equivalence relations (11)
$$Z_i^A(\tau )=Z_j^A(\tau )\mathrm{\Lambda }_i^j(\tau ),\psi _a^{}(\tau )=\lambda (\tau )\psi _a(\tau )$$
(22)
The global Poincaré symmetry leads to the on-shell conservation of the Hamiltonian counterparts of Poincaré generators $`𝐏_A,𝐌_{AB}`$
$$\begin{array}{c}𝐏_A=\overline{p}_A,𝐌^{AB}=x^A\overline{p}^Bx^B\overline{p}^A+S^{AB},\\ \\ S^{AB}=2i\{Z_i^A^{iB}\overline{Z}_i^A\overline{}^{iB}(\mathrm{\Sigma }^{AB})_b^a(\psi _a^b+\overline{\psi }_a\overline{}^b)\}\mathrm{\Phi }\end{array}$$
(23)
$`(\mathrm{\Sigma }_{AB})_a^b=\frac{1}{4}[\mathrm{\Gamma }_A,\mathrm{\Gamma }_B]_a^b`$ being the Lorentz generators in the spinor representation.
Let us now specify the expressions (19, 20) to the case of $`d=4`$. Then the spinning sector is parametrized by one variable $`\psi `$ defined modulo multiplication by a complex nonzero constant and subject to the Dirac equation. The expression for $`p_A`$ reads
$$p_A=\overline{p}_A(\psi ,\psi ^{})m\frac{(\stackrel{~}{\psi }\mathrm{\Gamma }_{AB}\psi )(\stackrel{~}{\psi }\mathrm{\Gamma }^B\psi )^{}+c.c.}{2\left|(\stackrel{~}{\psi }\mathrm{\Gamma }_C\psi )\right|^2}$$
(24)
and the consistency conditions take the form
$$\overline{p}^A(\stackrel{~}{\psi }\mathrm{\Gamma }_A\psi )=0,(\stackrel{~}{\psi }\mathrm{\Gamma }_{AB}\psi )=\frac{2}{m}\overline{p}_{[A}(\stackrel{~}{\psi }\mathrm{\Gamma }_{B]}\psi )$$
(25)
In terms of two-component Weyl spinors, $`\psi ^t=(\xi _a,\overline{\eta }^{\dot{a}})`$, the conditions (25) are equivalent to the following one:
$$\mathrm{𝐈𝐦}(\xi _a\eta ^a)=0$$
(26)
The Lagrangian (21) is specified as
$$=m\frac{\dot{x}_A(\sigma _{a\dot{a}}^A)(\xi ^a\overline{\xi }^{\dot{a}}+\eta ^a\overline{\eta }^{\dot{a}})}{2(\xi _a\eta ^a)}+is\frac{(\dot{\xi }_a\eta ^a)(\dot{\overline{\xi }}_a\overline{\eta }^a)}{(\xi _a\eta ^a)},$$
(27)
where $`m`$ and $`s`$ stand for the mass and spin of the particle. As is seen, the Lagrangian is invriant under the local projective transformations
$$\xi \alpha \xi ,\eta \overline{\alpha }\eta \alpha 𝐂\backslash \{0\}$$
(28)
and becomes singular whenever denominator $`\xi _a\eta ^a`$ comes to zero. To remove this singularity we can put the partial Lorenz invariant gauge on $`\xi `$ and $`\eta `$, breaking the invariance under rescalings with a real $`\alpha `$ <sup>e</sup><sup>e</sup>eThe higher dimensional generalization of this gauge, removing singularity in expression for the momenta (20), is obvious – $`|(\stackrel{~}{\psi }\mathrm{\Gamma }_{C(r)}\psi )|^2=m^2`$.:
$$\xi _a\eta ^a=m,$$
(29)
In this gauge the Lagrangian (27) of d=4 spinning particle takes a quite simple form
$$=\frac{1}{2}\dot{x}_A(\sigma _{a\dot{a}}^A)(\xi ^a\overline{\xi }^{\dot{a}}+\eta ^a\overline{\eta }^{\dot{a}})+\frac{2s}{m}\mathrm{𝐈𝐦}(\xi _a\dot{\eta }^a),$$
(30)
where the spinors $`\xi `$ and $`\eta `$ are assumed to be subjected to the one (comlex) holonomic constraint (29). The canonical momenta of the particle resulting from the Lagrangian looks like
$$p^A=\frac{1}{2}\sigma _{a\dot{a}}^A(\xi ^a\overline{\xi }^{\dot{a}}+\eta ^a\overline{\eta }^{\dot{a}})$$
(31)
and automatically satisfies mass-shell condition $`p^2=m^2`$ in view of (29). The representation (31) for the momenta of the d=4 massive spinning particle in term of two constrained Weyl spinors was originally considered in .
Now let us turn back to the original formulation (17) with unresolved momenta and consider the minimal coupling of the particle to an arbitrary background of gravitational and electromagnetic fields. For this end we introduce gauge fields of the vielbein $`e_\mu ^A`$ and torsion-free spin connection $`\omega _{\mu AB}`$ associated to the gravity and the electromagnetic potential $`A_\mu `$. Then the minimal covariantization of (17) reads
$$\begin{array}{c}S=\stackrel{~}{\theta },\\ \\ \stackrel{~}{\theta }=(p_Ae_\mu ^AeA_\mu )dx^\mu +D\mathrm{\Phi },\end{array}$$
(32)
where $`D`$ is the Lorentz covariant differential along the particle world-line
$$\begin{array}{c}DZ_i^A=dZ_i^A+dx^\mu \omega _\mu {}_{}{}^{A}{}_{B}{}^{}Z_{i}^{B},\\ \\ D\psi _a=d\psi _a+dx^\mu \omega _{\mu AB}\mathrm{\Sigma }_a^{ABb}\psi _b,\end{array}$$
(33)
and $`e`$ is the electric charge. The relations (2, 12) on the momentum $`p_A`$ are still assumed to hold. The action (32) generates the following equations of motion
$$\begin{array}{c}\dot{x}^\nu =\mu e_A^\nu p^A,\frac{Dp_A}{d\tau }+\mu eF_{AB}p^B=(\mu /4)R_{ABCD}p^BS^{CD},\\ \\ \frac{DZ_i^A}{d\tau }=0,\frac{D\psi _a}{d\tau }=0\end{array}$$
(34)
Here $`F_{\mu \nu }`$ is the strength tensor of the electromagnetic field and $`R_{\alpha \beta \gamma \delta }`$ is a curvature of the space-time. Thus the dynamics of the particle in the curved space-time is described by first two equations while the motion in the spinning sector reduces to the parallel transport along the world-line.
To conclude, let us note that the gauge symmetry (22), being required to survive on the quantum level, implies the restriction on the possible values of the parameters $`s_i`$ entering the Kähler potential. Proceeding by analogy with the integer spin case one can show that in quantum theory $`s_i,i=1,\mathrm{}r1`$ are constrained to be integer and $`s_r`$ \- (half-)integer numbers.
## 3 Quantization
In this section we will present the covariant quantization of the model described. Let us start with the Lagrangian (21) complemented with the conditions (13) and (19). This Lagrangian leads to the following primary constraints
$$\begin{array}{c}T_A=p_A\overline{p}_A0,\\ \\ _A^i=q_A^i+i_A^i\mathrm{\Phi }0,\overline{}_A^i=(_A^i)^{},\\ \\ ^a=\pi ^a+i^a\mathrm{\Phi }0,\overline{}^{\dot{a}}=(^a)^{}\end{array}$$
(35)
(as well as the above mentioned holonomic ones). Here $`p_A`$, $`q_A^i`$ and $`\pi ^a`$ are the momenta conjugated to $`x^A`$, $`Z_i^A`$ and $`\psi _a`$ respectively, $`\overline{p}_A`$ is defined by rels. (20). The system (35) allows a transition to the more suitable basis of constraints where relations (19) have been already accounted
$$\begin{array}{c}T_i=(p,Z_i)0,T_a=p_A\mathrm{\Gamma }^A\psi m\psi 0,\\ \\ \overline{}_A^i0,\overline{}^{\dot{a}}0\end{array}$$
(36)
and the complex conjugated constraints are also implied. Stabilization of the constraints yields the mass-shell condition
$$T=p_Ap^A+m^20$$
(37)
It should be noted that the constraints (36, 37) are not independent. Moreover, only first class constraints amongst (36)<sup>f</sup><sup>f</sup>f Of course, (37) is also first class. can be covariantly extracted:
$$\begin{array}{c}\mathrm{\Pi }_j^i=(Z_j,q^i)+i\delta _j^il_i,\overline{\mathrm{\Pi }}_j^i=(\mathrm{\Pi }_j^i)^{},ij,\\ \\ \mathrm{\Pi }=(\psi _a\pi ^a)+il_r,\overline{\mathrm{\Pi }}=(\mathrm{\Pi })^{},\\ \\ l_i=\underset{k=i}{\overset{r1}{}}s_k,l_r=2s_r\end{array}$$
(38)
Constraints $`\mathrm{\Pi }_j^i`$ and $`\mathrm{\Pi }`$ generate the gauge transformations (11). Following the covariant quantization scheme, the set of variables $`(x,p,Z,q,\psi ,\pi )`$ is associated to the self-adjoint operators acting in a Hilbert space of the particle states. The physical states $`|\mathrm{\Psi }`$ are singled out from the space of smooth functions on $`𝐑^{d1,1}\times 𝐂^{(r1)d}\times 𝐂^{2^{d/2}}`$ by imposing the first class constraint operators and a half of the second class ones. Since the latter cannot be presented explicitly we will impose on the physical states the operatorial counterparts of the expressions (36) (accounting thereby the half of the second class constraints), together with $`\widehat{\overline{\mathrm{\Pi }}}_j^i`$ and $`\widehat{\overline{\mathrm{\Pi }}}`$. As a result, the physical states are annihilated by the operators $`\widehat{T}`$, $`\widehat{T}_a`$, $`\widehat{T}_i`$ , $`\widehat{\mathrm{\Pi }}`$, $`\widehat{\mathrm{\Pi }}_j^i`$, $`\widehat{\overline{}}^{\dot{a}}`$, $`\widehat{\overline{}}_A^i`$. In the coordinate representation for $`(Z,q,\psi ,\pi )`$ and the momentum one for $`(x,p)`$
$$q_A^ii_A^i,\pi ^ai^a,x^Ai^A$$
(39)
the constraint equations for the wave function $`\mathrm{\Psi }(Z,\overline{Z},\psi ,\psi ^{},p)`$ take the following explicit form:
$$(Z_j^A_A^i\delta _j^in_i)\mathrm{\Psi }=0,(\psi _a^an)\mathrm{\Psi }=0$$
(40)
$$\overline{}_A^i\mathrm{\Psi }+(\overline{}_A^i\mathrm{\Phi })\mathrm{\Psi }=0,\overline{}^{\dot{a}}\mathrm{\Psi }+(\overline{}^{\dot{a}}\mathrm{\Phi })\mathrm{\Psi }=0$$
(41)
The wave function is defined on the surface (12, 13). Note that the constants $`n`$, $`n_i`$ may differ from their classical values $`l`$, $`l_i`$ due to the different operator ordering prescriptions for $`\widehat{\mathrm{\Pi }}`$ and $`\widehat{\mathrm{\Pi }}_j^i`$. We fix the ambiguity in the factor ordering by the requirement that the physical wave functions should remain unchanged under the gauge transformations (11) which implies vanishing of $`n`$ and $`n_i`$. Any other ordering will lead to the unitary equivalent theory. The general solution for (41) reads
$$\mathrm{\Psi }(Z,\overline{Z},\psi ,\psi ^{},p)=\mathrm{exp}(\mathrm{\Phi }(Z,\overline{Z},\psi ,\psi ^{}))\mathrm{\Theta }(Z,\psi ,p)$$
(42)
Substituting (42) in (40) one gets
$$(Z_j^A_A^i\delta _j^il_i)\mathrm{\Theta }=0,(\psi _a^al_r)\mathrm{\Theta }=0$$
(43)
Since the manifold $`𝒪_𝐬`$ is compact the eigenvalues $`l_i`$, $`l_r`$ prove to be integers with eigenfunctions $`\mathrm{\Theta }`$ represented by the polynomials of the form
$$\mathrm{\Theta }(Z,\psi ,p)=\mathrm{\Theta }(p)_{A(l_1)B(l_2)\mathrm{}C(l_{r1})a(l_r)}Z_1^{A(l_1)}Z_2^{B(l_2)}\mathrm{}Z_{r1}^{C(l_{r1})}\stackrel{~}{\psi }^{a(l_r)},$$
(44)
where we denote $`Z_1^{A(l)}Z_1^{A_1}\mathrm{}Z_1^{A_l}`$, $`\stackrel{~}{\psi }{}_{}{}^{a(l)}\stackrel{~}{\psi }{}_{}{}^{a_1}\mathrm{}\stackrel{~}{\psi }^{a_l}`$ ; the symmetry of the indices is described by the Young diagram
Notice that, by virtue of relations (11, 12, 13), the coefficients $`\mathrm{\Theta }(p)_{A\mathrm{}Ca\mathrm{}}`$ are assumed to be $`p`$-transversal, traceless and $`\mathrm{\Gamma }`$-traceless
$$p^A\mathrm{\Theta }(p)_\mathrm{}A\mathrm{}=0,\eta ^{AB}\mathrm{\Theta }(p)_{\mathrm{}A\mathrm{}B\mathrm{}}=0,\mathrm{\Gamma }_a^{Ab}\mathrm{\Theta }(p)_{\mathrm{}A\mathrm{}b\mathrm{}}=0$$
(45)
In spinor indices $`\mathrm{\Theta }`$ is subject to the Dirac equation
$$p_A\mathrm{\Gamma }_a^{Ab}\mathrm{\Theta }(p)_\mathrm{}b\mathrm{}=m\mathrm{\Theta }(p)_\mathrm{}a\mathrm{}$$
(46)
Equations (45, 46) constitute together the full set of $`d`$ -dimensional relativistic wave equations on irreducible massive spin tensor fields which are obtained after Fouriau transform of $`\mathrm{\Theta }(p)`$. The space $`_{m,𝐬}`$ of functions $`\mathrm{\Psi }`$ (42) representing the particle states is endowed with the Hilbert space structure with respect to the following Hermitian inner product
$$<\mathrm{\Psi }_1|\mathrm{\Psi }_2>=\underset{p^2=m^2}{}\frac{d𝐩}{p_0}\underset{𝒪_𝐬}{}𝑑\mu \overline{\mathrm{\Psi }}_1\mathrm{\Psi }_2$$
(47)
where
$$d\mu =dd\mathrm{\Phi }dd\mathrm{\Phi }\mathrm{}dd\mathrm{\Phi }$$
is the Liouville measure on $`𝒪_𝐬`$. Integration over the spinning degrees of freedom may be performed explicitly resulting with the standard field-theoretical inner product
$$<\mathrm{\Psi }_1|\mathrm{\Psi }_2>=N\underset{p^2=m^2}{}\frac{d𝐩}{p_0}\overline{\mathrm{\Theta }}(p)_1^{A(l_1)B(l_2)\mathrm{}C(l_{r1})a(l_r)}\mathrm{\Theta }(p)_{2A(l_1)B(l_2)\mathrm{}C(l_{r1})a(l_r)}$$
(48)
where $`\overline{\mathrm{\Theta }}^{\mathrm{}a(l_r)}=\mathrm{\Theta }^{\mathrm{}}{}_{\dot{a}(l_r)}{}^{}(\mathrm{\Gamma }_0)_{}^{\dot{a}_1a_1}\mathrm{}(\mathrm{\Gamma }_0)^{\dot{a}_{l_r}a_{l_r}}`$ and $`N`$ is a constant depending on spin of the particle.
Thus the covariant operatorial quantization of the model yields the irreducible representation of the Poincaré group with quantum numbers fixed by the constants originally entering the Kähler potential. We would like also to note that the procedure of geometric quantization provides another way of quantizing these systems with the same result (see where the model of integer spin massive particle is quantized within this approach).
## 4 Conclusion
Let us discuss some possible links between the spinning particle model suggested and some problems of interest in other topics of the field theory and mention some open questions related to the paper results.
In this paper we have extended the description of the integer spin particle to a general (half-)integer case in arbitrary dimension. The key technical step of our construction is the new covariant vector-spinor parametrization for the phase space of spin $`𝒪_𝐬`$, allowing to obtain, upon covariant quantization, arbitrary half-ineger spin representations along with the integer ones. This construction may be thought of as a minimal ”spinorization” of the previous one in a sence that only one vector is replaced by a spinor. In principle, it is possible to trade all the vector coordinates for the spinor variables subject to a certain set of constraints, which may seem to be more fundamental for dealing with the half-integer spin representations. Then the harmonic expansion of physical wave functions on these variables would define (upon quantization) the irreducible spin-tensor fields similarly to (44). To the best of our knowledge, the respective analysis of irreducibility for the general spin-tensor fields is yet unknown in the higher dimensions, what considerably hinders the determination of an appropriate set of constraints on spinor variables.
As to the problem of interaction between particle and axternal fields, we may mention that, at least at the formal level, it is solved in this paper for the minimal coupling to a general gravity and electromagnetic background fields. The extension can be immediately got to the particle interaction with the dynamical fields by adding the free field Lagrangian in the action. However, the less formal aspect related to the self-accelerating problem remains unclear now in the selfconsistent theory of the particle coupled to the dynamical fields. In ref. has been shown that the selfacceleration of the particle does not occur when it is coupled to a very special spectrum of the fields. This special field spectrum is known for a $`d=4`$ spinless particle only. For higher dimensions and/or nonzero spin, the solution to the selfacceleration problem could be different. It should be mentioned also that the interaction to the non-abelian gauge fields requires to equip the particle with isospinning degrees of freedom. This could be done, in general, along the same lines as it is performed for a genuine spin in this paper, although the isospin requires different inner manifold instead of $`𝒪_𝐬`$. In this way, the isospin-shell conditions should appear as phase space constraints. However, it is unclear from the outset, whether the interplay between spin and isospin in the constraint algebra would be consistent to an arbitrary Yang-Mills background field.
Another problem, where this model may seem to gain some importance, is the relationship between strings and spinning particles. This relationship is commonly known at the level of the quantum string state spectrum which includes an infinite number of massive excitations of various spins. These excitations are usually thought about as states of certain spinning particles. However, it remains yet unclear how the nonzero string modes (subject to the Virasoro constraints) may form the spinning sector of the particle phase space. In other words, the question is how the surface of Virasoro constraints in the string phase space is stratified into the spinning particle presymplectic manifolds $``$. Comprehension of the structure of this stratification seems to be relevant for study of a reduction procedure in string theory. The geometry of the spinning particle phase space, revealed in this paper, forms a basis to study the relationship between strings and particles in this direction.
## 5 Acknowledgments
This work is partially supported by the grant Joint INTAS-RFBR 95-829 and RFBR 98-02-16261. A. A. Sharapov appreciates the financial support from the INTAS under the grant YSF 98-153. |
warning/0002/nucl-th0002008.html | ar5iv | text | # The Quasi-Molecular Stage of Ternary Fission
## 1 Introduction
The light particle accompanied fission was discovered in 1946, when the track of a long-range particle (identified by Farwell et al. to be <sup>4</sup>He) almost perpendicular to the short tracks of heavy and light fragments was observed in a photographic plate. The fission was induced by bombarding <sup>235</sup>U with slow neutrons from a Be target at a cyclotron. The largest yield in such a rare process (less than one event per 500 binary splittings) is measured for $`\alpha `$-particles, but many other light nuclei (from protons to oxygen or even calcium isotopes) have been identified in both induced- and spontaneous fission phenomena. If $`A_1`$ and $`A_2`$ are the mass numbers of the heavy fragments (assume $`A_1A_2`$), then usually the mass of the light particle $`A_3<<A_2`$. The “true” ternary fission, in which $`A_1A_2A_3`$, has not yet been experimentally detected.
Many properties of the binary fission process have been explained within the liquid drop model (LDM); others like the asymmetric mass distribution of fragments and the ground state deformations of many nuclei, could be understood only after adding the contribution of shell effects. As it was repeatedly stressed (see and the references therein), shell effects proved also to be of vital importance for cluster radioactivities predicted in 1980.
The total kinetic energy (TKE) of the fragments, in the most frequently detected binary or ternary fission mechanism, is smaller than the released energy ($`Q`$) by about 25–35 MeV, which is used to produce deformed and excited fragments. These then emitt neutrons (each with a binding energy of about 6 MeV) and $`\gamma `$-rays. From time to time a “cold” fission mechanism is detected, in which the TKE exhausts the $`Q`$-value, hence no neutrons are emitted, and the fragments are produced in or near their ground-state. The first experimental evidence for cold binary fission in which its TKE exhaust $`Q`$ was reported in 1981. Larger yields were measured in trans-Fm ($`Z100`$) isotopes, where the phenomenon was called bimodal fission.
The correlated fragment pairs in cold ternary ($`\alpha `$\- and <sup>10</sup>Be accompanied spontaneous fission of <sup>252</sup>Cf) processes were only recently discovered, by measuring triple $`\gamma `$ coincidences in a modern large array of $`\gamma `$-ray detectors (GAMMASPHERE). The fragments are identified by their $`\gamma `$-ray spectra. Among other new aspects of the fission process seen for the first time with this new technique, one should mention the double fine structure, and the triple fine structure in binary and ternary fission.
A particularly interesting feature, observed both in <sup>10</sup>Be- and <sup>12</sup>C accompanied cold fission of <sup>252</sup>Cf is related to the width of the light particle $`\gamma `$-ray spectrum (see Figs. 1 and 2). For example, the 3.368 MeV $`\gamma `$ line of <sup>10</sup>Be, with a lifetime of 125 fs is not Doppler-broadened, as it should be if it would be emitted when <sup>10</sup>Be is in flight (taking about 1 ns to reach the detector). A plausible suggestion was made, that the absence of Doppler broadening is related to a trapping of <sup>10</sup>Be in a potential well of nuclear molecular character.
Quasi-molecular configurations of two nuclei have been suggested as a natural explanation for the resonances measured in <sup>12</sup>C$`+^{12}`$C scattering and reactions. There are also other kinds of such binary molecules (see and references therein), like spontaneously fissioning shape-isomers. The above mentioned experiments can be considered as the first evidence for a more complex quasi-molecular configuration of three nuclei. The purpose of the present lecture is to show, within a phenomenological three-center model, that a minimum which could explain the existence of these quasi-molecules is produced in the potential barrier, when the formation of the light particle occurs in the neck between the two heavier fragments. In this way we extend to ternary fission our unified approach of cold fission, cluster radioactivities, and $`\alpha `$-decay.
## 2 Shape Parametrization
The shape parametrization with one deformation parameter as follows has been suggested from the analysis of different aligned and compact configurations of fragments in touch. A lower potential barrier for the aligned cylindrically-symmetric shapes with the light particle between the two heavy fragments, is a clear indication that during the deformation from an initial parent nucleus to three final nuclei, one should arrive at such a scission point. In order to reach this stage we shall increase continuously the separation distance, $`R`$, between the heavy fragments, while the radii of the heavy fragment and of the light particle are kept constant, $`R_1=`$ constant, $`R_3=`$ constant. Unlike in the previous work, we now adopt the following convention: $`A_1A_2A_3`$. The hadron numbers are conserved: $`A_1+A_2+A_3=A`$.
At the beginning (the neck radius $`\rho _{neck}R_3`$) one has a two-center evolution (see Fig. 3) until the neck between the fragments becomes equal to the radius of the emitted particle, $`\rho _{neck}=\rho (z_{s1})_{R=R_{ov3}}=R_3`$. This Eq. defines $`R_{ov3}`$ as the separation distance at which the neck radius is equal to $`R_3`$. By placing the origin in the center of the large sphere, the surface equation in cylindrical coordinates is given by:
$$\rho _s^2=\{\begin{array}{ccc}& R_1^2z^2\hfill & ,R_1zz_{s1}\hfill \\ & R_2^2(zR)^2\hfill & ,z_{s1}zR+R_2\hfill \end{array}$$
(1)
Then for $`R>R_{ov3}`$ the three center starts developing by decreasing progressively with the same amount the two tip distances $`h_1+h_{31}=h_{32}+h_2`$. Besides this constraint, one has as in the binary stage, volume conservation and matching conditions. The $`R_2`$ and the other geometrical quantities are determined by solving numerically the corresponding system of algebraic equations. By assuming spherical nuclei, the radii are given by $`R_j=1.2249A_j^{1/3}`$ fm ($`j=0,1,3`$), $`R_{2f}=1.2249A_2^{1/3}`$ with a radius constant $`r_0=1.2249`$ fm, from Myers-Swiatecki’s variant of LDM. Now the surface equation can be written as
$$\rho _s^2=\{\begin{array}{ccc}& R_1^2z^2\hfill & ,R_1zz_{s1}\hfill \\ & R_3^2(zz_3)^2\hfill & ,z_{s1}zz_{s2}\hfill \\ & R_2^2(zR)^2\hfill & ,z_{s2}zR+R_2\hfill \end{array}$$
(2)
and the corresponding shape has two necks and two separating planes. Some of the important values of the deformation parameter $`R`$ are the initial distance $`R_i=R_0R_1`$, and the touching-point one, $`R_t=R_1+2R_3+R_{2f}`$. There is also $`R_{ov3}`$, defined above, which allows one to distinguish between the binary and ternary stage.
## 3 Deformation Energy
According to the LDM, by requesting zero energy for a spherical shape, the deformation energy, $`E^u(R)E^0`$, is expressed as a sum of the surface and Coulomb terms
$$E_{def}^u(R)=E_s^0[B_s(R)1]+E_C^0[B_C(R)1]$$
(3)
where the exponent <sup>u</sup> stands for uniform (fragments with the same charge density as the parent nucleus), and <sup>0</sup> refers to the initial spherical parent. In order to simplify the calculations, we initially assume the same charge density $`\rho _{1e}=\rho _{2e}=\rho _{3e}=\rho _{0e}`$, and at the end we add the corresponding corrections. In this way we perform one numerical quadrature instead of six. For a spherical shape $`E_s^0=a_s(1\kappa I^2)A^{2/3}`$ ; $`I=(NZ)/A`$; $`E_C^0=a_cZ^2A^{1/3}`$, where the numerical constants of the LDM are: $`a_s=17.9439`$ MeV, $`\kappa =1.7826`$, $`a_c=3e^2/(5r_0)`$, $`e^2=1.44`$ MeV$``$fm.
The shape-dependent, dimensionless surface term is proportional to the surface area:
$$B_s=\frac{E_s}{E_s^0}=\frac{d^2}{2}\underset{1}{\overset{+1}{}}\left[y^2+\frac{1}{4}\left(\frac{dy^2}{dx}\right)^2\right]^{1/2}𝑑x$$
(4)
where $`y=y(x)`$ is the surface equation in cylindrical coordinates with -1, +1 intercepts on the symmetry axis, and $`d=(z^{\prime \prime }z^{})/2R_0`$ is the seminuclear length in units of $`R_0`$. Similarly, for the Coulomb energy one has
$$B_c=\frac{5d^5}{8\pi }\underset{1}{\overset{+1}{}}𝑑x\underset{1}{\overset{+1}{}}𝑑x^{}F(x,x^{})$$
(5)
$`F(x,x^{})`$ $`=`$ $`\{yy_1[(K2D)/3]`$ (6)
$`\left[2(y^2+y_1^2)(xx^{})^2+{\displaystyle \frac{3}{2}}(xx^{})\left({\displaystyle \frac{dy_1^2}{dx^{}}}{\displaystyle \frac{dy^2}{dx}}\right)\right]+`$
$`K\{y^2y_1^2/3+[y^2{\displaystyle \frac{xx^{}}{2}}{\displaystyle \frac{dy^2}{dx}}][y_1^2{\displaystyle \frac{xx^{}}{2}}{\displaystyle \frac{dy_1^2}{dx^{}}}]\}\}a_\rho ^1`$
$`K`$, $`K^{}`$ are the complete elliptic integrals of the 1st and 2nd kind
$$K(k)=\underset{0}{\overset{\pi /2}{}}(1k^2\mathrm{sin}^2t)^{1/2}𝑑t;K^{}(k)=\underset{0}{\overset{\pi /2}{}}(1k^2\mathrm{sin}^2t)^{1/2}𝑑t$$
(7)
and $`a_\rho ^2=(y+y_1)^2+(xx^{})^2`$, $`k^2=4yy_1/a_\rho ^2`$, $`D=(KK^{})/k^2`$.
The new minimum, which can be seen in Fig. 4 at a separation distance $`R=R_{mint}>R_{ov3}`$, is the result of a competition between the Coulomb- and surface energies. At the beginning ($`R<R_{mint}`$) the Coulomb term is stronger, leading to a decrease in energy, but later on ($`R>R_{mint}`$) the light particle formed in the neck posses a surface area increasing rapidly, so there is also an increase in energy up to $`R=R_t`$.
Now let us analyse the influence of various corrections, which could in principle alter this image. After performing numerically the integrations, we add the following corrections: for the difference in charge densities reproducing the touching point values; for experimental masses reproducing the $`Q_{exp}`$-value at $`R=R_i`$, when the origin of energy corresponds to infinite separation distances between fragments, and the phenomenological shell corrections $`\delta E`$
$$E_{LD}(R)=E_{def}^u(R)+(Q_{th}Q_{exp})f_c(R)$$
(8)
where $`f_c(R)=(RR_i)/(R_tR_i)`$, and
$$Q_{th}=E_s^0+E_C^0\underset{1}{\overset{3}{}}(E_{si}^0+E_{Ci}^0)+\delta E^0\underset{1}{\overset{3}{}}\delta E^i$$
(9)
The correction increases gradually (see Fig. 4 and Fig. 5) with $`R`$ up to $`R_t`$ and then remains constant for $`R>R_t`$. The barrier height increases if $`Q_{exp}<Q_{th}`$ and decreases if $`Q_{exp}>Q_{th}`$. In this way, when one, two, or all final nuclei have magic numbers of nucleons, $`Q_{exp}`$ is large and the fission barrier has a lower height, leading to an increased yield. In a binary decay mode like cluster radioactivity and cold fission, this condition is fulfilled when the daughter nucleus is <sup>208</sup>Pb and <sup>132</sup>Sn, respectively.
## 4 Shell Corrections and Half-lives
Finally we also add the shell terms
$$E(R)=E_{LD}(R)+\delta E(R)\delta E^0$$
(10)
Presently there is not available any microscopic three-center shell model reliably working for a long range of mass asymmetries. This is why we use a phenomenological model, instead of the Strutinsky’s method, to calculate the shell corrections. The model is adapted after Myers and Swiatecki. At a given $`R`$, we calculate the volumes of fragments and the corresponding numbers of nucleons $`Z_i(R),N_i(R)`$ ($`i=1,2,3`$), proportional to the volume of each fragment. Fig. 3 illustrates the evolution of shapes and of the fragment volumes. Then we add for each fragment the contribution of protons and neutrons
$$\delta E(R)=\underset{i}{}\delta E_i(R)=\underset{i}{}[\delta E_{pi}(R)+\delta E_{ni}(R)]$$
(11)
which are given by
$$\delta E_{pi}=Cs(Z_i);\delta E_{ni}=Cs(N_i)$$
(12)
where
$$s(Z)=F(Z)/[(Z)^{2/3}]cZ^{1/3}$$
(13)
$$F(n)=\frac{3}{5}\left[\frac{N_i^{5/3}N_{i1}^{5/3}}{N_iN_{i1}}(nN_{i1})n^{5/3}+N_{i1}^{5/3}\right]$$
(14)
in which $`n(N_{i1},N_i)`$ is either a current $`Z`$ or $`N`$ number and $`N_{i1},N_i`$ are the closest magic numbers. The constants $`c=0.2`$, $`C=6.2`$ MeV were determined by fit to the experimental masses and deformations. The variation with $`R`$ is calculated as
$$\delta E(R)=\frac{C}{2}\left\{\underset{i}{}[s(N_i)+s(Z_i)]\frac{L_i(R)}{R_i}\right\}$$
(15)
where $`L_i(R)`$ are the lengths of the fragments along the axis of symmetry, at a given separation distance $`R`$. During the deformation, the variation of separation distance between centers, $`R`$, induces the variation of the geometrical quantities and of the corresponding nucleon numbers. Each time a proton or neutron number reaches a magic value, the correction energy passes through a minimum, and it has a maximum at midshell (see Fig. 5 and Fig. 6). The first narrow minimum appearing in the shell correction energy $`\delta E`$ in Fig. 5, at $`R=R_{min1b}2.6`$ fm, is the result of almost simultaneously reaching the magic numbers $`Z_1=20`$, $`N_1=28`$, and $`Z_2=82`$, $`N_2=126`$. The second, more
shallower one around $`R=R_{min2b}7.2`$ fm corresponds to a larger range of $`R`$-values for which $`Z_1=50`$, $`N_1=82`$, $`Z_2=50`$, $`N_2=82`$ are not obtained in the same time. In the region of the new minimum, $`R=R_{mint}`$, for light-particle accompanied fission, the variation of the shell correction energy is very small, hence it has no major consequence. One can say that the quasimolecular minimum is related to the collective properties (liquid-drop like behavior). On the other side, for “true” ternary process (see the bottom part of Fig. 6) both minima appear in this range of values, but no such LDM effect was found there. In order to compute the half-life of the quasi-molecular state, we have first to search for the minimum $`E_{min}`$ in the quasimolecular well, from $`R_{ov3}`$ to $`R_t`$, and then to add a zero point vibration energy, $`E_v`$: $`E_{qs}=E_{min}+E_v`$.
The half-life, $`T`$, is expressed in terms of the barrier penetrability, $`P`$, which is calculated from an action integral, $`K`$, given by the quasi-classical WKB approximation
$$T=\frac{h\mathrm{ln}2}{2E_vP};P=exp(K)$$
(16)
where $`h`$ is the Planck constant, and
$$K=\frac{2}{\mathrm{}}\underset{R_a}{\overset{R_b}{}}\sqrt{2\mu [E(R)E_{qs}]}𝑑R$$
(17)
in which $`R_a,R_b`$ are the turning points, defined by $`E(R_a)=E(R_b)=E_{qs}`$ and the nuclear inertia is roughly approximated by the reduced mass $`\mu =m[(A_1A_2+A_3A)/(A_1+A_2)]`$, where $`m`$ is the nucleon mass, $`\mathrm{log}[(h\mathrm{ln}2)/2]=20.8436`$, $`\mathrm{log}e=0.43429`$ and $`\sqrt{8m/\mathrm{}^2}=0.4392`$ MeV$`{}_{}{}^{1/2}\times `$fm<sup>-1</sup>.
The results of our estimations for the half-lives of some quasimolecular states formed in the <sup>10</sup>Be- and <sup>12</sup>C accompanied fission of <sup>252</sup>Cf are given in Table 1. They are of the order of 1 ns and 1 ms, respectively, if we ignore the results for a division with heavy fragment <sup>132</sup>Sn, which was not measured due to very high first excited state. Consequently the new minimum we found can qualitatively explain the quasimolecular nature of the narrow line of the <sup>10</sup>Be $`\gamma `$-rays.
It is interesting to note that the trend toward a split into two, three, or four nuclei (the lighter ones formed in a long neck between the heavier fragments) has been theoretically demonstrated by Hill, who investigated the classical dynamics of an incompressible, irrotational, uniformly charged liquid drop. No mass asymmetry was evidenced since any shell effect was ignored.
In conclusion, we should stress that a quasimolecular stage of a light-particle accompanied fission process, for a limited range of sizes of the three partners, can be qualitatively explained within the liquid drop model.
### Acknowledgments
We are grateful to M. Mutterer for enlightening discussions. |
warning/0002/hep-th0002188.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Recently, Maldacena proposed a duality between type IIB string theory and (super-) conformal field theory (CFT) on the boundary of anti-de Sitter space ($`AdS`$) , further elaborated on in Refs. . This remarkable conjecture has induced a tremendous activity in theoretical high energy physics. In the cases involving $`AdS_3`$, see e.g. Refs. and the review , the corresponding space-time CFT is two-dimensional and thus possesses many well known properties.
The boundary of $`AdS_3`$ enjoys conformal symmetry, as recognized by Brown and Henneaux . In the work , Giveon et al have constructed explicitly the generators of the space-time conformal algebra from the string theory on $`AdS_3`$. The construction starts from the world sheet $`SL(2)`$ current algebra. In terms of the Wakimoto free field realization of that, they have provided a general expression for the Virasoro generators and computed the central charge.
In Ref. Ito has succeeded in constructing superconformal algebras (SCAs) on the boundary of $`AdS_3`$ and again the construction starts from a world sheet current algebra. However, in order to obtain an extended conformal symmetry in space-time one needs to consider higher world sheet current (super-)algebras than $`SL(2)`$. Thus, Ito has found that the appropriate Lie superalgebras leading to $`N=1`$, 2 and 4 superconformal algebras are $`osp(1|2)`$, $`sl(2|1)`$ and $`sl(2|2)`$, respectively. The explicit constructions are based on generalized Wakimoto free field realizations of the associated affine current superalgebras . A related approach to construct $`N=1`$, 2 and 4 SCAs is discussed in in which also one- and two-point functions and some unitary representations are considered. In Ref. space-time $`N=3`$ superconformal theories are studied.
The objective of the present paper is to take a first step in the direction of classifying the SCAs that may be induced by string theory on $`AdS_3`$, thus generalizing the work by Ito . A main property of an appropriate world sheet current (super-)algebra is that the bosonic part may be decomposed as $`G=SL(2)G^{}`$. This is necessary as we want the free fields in the Wakimoto realization of the embedded $`SL(2)`$ to be considered as coordinates on $`AdS_3`$. As discussed in Ref. , a purely bosonic world sheet algebra with such a decomposition leads immediately to an affine Lie algebra on the boundary of $`AdS_3`$. The present paper is devoted to discussing the class of superconformal algebras that may be constructed starting from affine $`SL(2|N/2)`$ current superalgebras having as bosonic part $`SL(2)SL(N/2)U(1)`$. By construction, $`N`$ is even.
We construct explicitly the Virasoro generators, the $`N`$ supercurrents, and the generators of an internal $`SL(N/2)U(1)`$ Kac-Moody algebra. For $`N=2`$, the $`SL(N/2)`$ and $`U(1)`$ current algebras collapse to a single $`U(1)`$ current algebra. This conventional $`N=2`$ superconformal algebra has already been obtained by Ito . For higher $`N`$ we turn to a classical limit in which the generators may be substituted by first order linear differential operators. The resulting classical SCAs are center-less and are shown to include classes of primary generators of half-integer weights which are all smaller than 2. Some weights are negative for $`N6`$. SCAs based on free field realizations and with generically non-vanishing central charges are denoted quantum SCAs as opposed to such classical SCAs. Thus, our use of the notion quantum is not in the quantum group sense of $`q`$-deformations.
A new and important property of the construction is that for $`N2`$ it treats the supercurrents asymmetrically. This is illustrated in the case $`N=4`$ where the classical SCA is completed and found to be of a new and asymmetric form. Thus, it is not included in the standard classification of $`N=4`$ SCAs . In particular, it deviates essentially from the small $`N=4`$ SCA which has otherwise been announced to be the result of a construction similar to the one employed in the present paper. This is argued not to be the correct result. The full quantum $`N=4`$ SCA with generic central charge will be presented elsewhere . There we shall also show that the small $`N=4`$ SCA may be obtained by replacing the original world sheet $`SL(2|2)`$ current superalgebra by the related $`SL(2|2)/U(1)`$ current superalgebra.
A complete classification along the lines indicated is reached when the SCAs induced by any world sheet current superalgebra with $`SL(2)G^{}`$ decomposable bosonic part have been constructed. We anticipate that the techniques employed in the present paper may be enhanced to cover the general case and hope to come back elsewhere with a discussion on this generalization.
As pointed out in Ref. , BRST invariance of the construction of the space-time conformal algebra is equivalent to requiring the Virasoro generators to be primary fields of weight one with respect to the world sheet energy-momentum tensor, ensuring that the integrated fields commute with the world sheet Virasoro algebra. This carries over to the superconformal case, and we shall verify that the generators of our SCAs meet the requirement of being primary of weight one with respect to the world sheet current superalgebra Sugawara tensor.
The algebras constructed in Ref. are simpler than the ones by Ito as they are based on smaller Lie superalgebras. For example, the $`N=4`$ SCA is constructed from an $`sl(2|1)`$ Lie superalgebra. However, the central charges are essentially fixed, and the algebras are in general not ensured to be BRST invariant<sup>1</sup><sup>1</sup>1We thank O. Andreev for pointing out that the $`N=4`$ SCA in Ref. is nevertheless BRST invariant..
The remaining part of this paper is organized as follows. In Section 2 we review the construction of the Virasoro algebra and the immediate extension to an affine Lie algebra .
In Section 3 we introduce our notation for Lie superalgebras and their associated current superalgebras, and review the free field realizations of the latter obtained in Ref. .
In Section 4 we provide our explicit construction of the supercurrents, the Virasoro generators, and the generators of the internal $`SL(N/2)U(1)`$ Kac-Moody algebra.
In Section 5 BRST invariance of the construction is addressed.
In Section 6 we discuss the classical $`N`$ extended SCAs and write down the explicit result for $`N=4`$.
Section 7 contains concluding remarks, while details on the Lie superalgebra $`sl(2|M)`$ are given in Appendix A.
## 2 Virasoro Algebra
The standard Wakimoto free field realization of the affine $`SL(2)`$ current algebra with level $`k^{}`$ is
$`E`$ $`=`$ $`\beta `$
$`H`$ $`=`$ $`2\gamma \beta +\sqrt{k^{}+2}\phi `$
$`F`$ $`=`$ $`\gamma ^2\beta +\sqrt{k^{}+2}\gamma \phi +k^{}\gamma `$ (1)
Here and throughout the paper, normal ordering is implicit. The operator product expansions (OPEs) of the ghost fields $`\beta ,\gamma `$ and the bosonic scalar field $`\phi `$ are
$$\beta (z)\gamma (w)=\frac{1}{zw},\phi (z)\phi (w)=2\mathrm{ln}(zw)$$
(2)
where regular terms have been omitted. In Ref. it is shown that the world sheet $`SL(2)`$ current algebra with level $`k^{}`$ induces the Virasoro algebra
$$[L_n,L_m]=(nm)L_{n+m}+\frac{c}{12}(n^3n)\delta _{n+m,0}$$
(3)
on the boundary of $`AdS_3`$. The generators are given by
$`L_n`$ $`=`$ $`{\displaystyle \frac{dz}{2\pi i}_n(z)}`$
$`_n`$ $`=`$ $`a_+(n)\gamma ^{n+1}E+a_3(n)\gamma ^nH+a_{}(n)\gamma ^{n1}F`$ (4)
with constants<sup>2</sup><sup>2</sup>2Note that the conventions used here differ slightly from the ones used in Ref. , which is also reflected in a sign discrepancy in the central charge (6).
$`a_+(n)`$ $`=`$ $`{\displaystyle \frac{1}{2}}(1n)n`$
$`a_3(n)`$ $`=`$ $`{\displaystyle \frac{1}{2}}(1n)(1+n)`$
$`a_{}(n)`$ $`=`$ $`{\displaystyle \frac{1}{2}}n(1+n)`$ (5)
The central charge is found to be
$$c=6k^{}p$$
(6)
where $`p`$ is the integer winding number
$$p=\frac{dz}{2\pi i}\frac{\gamma }{\gamma }$$
(7)
BRST invariance requires $`_n`$ to be conformal primary of weight 1 with respect to the world sheet Virasoro generator
$$T=\gamma \beta +\frac{1}{2}\phi \phi \frac{1}{2\sqrt{k^{}+2}}^2\phi $$
(8)
This property is readily verified.
The Virasoro algebra is immediately extended to an affine Lie algebra if the world sheet affine $`SL(2)`$ current algebra is replaced by an affine $`G=SL(2)G^{}`$ current algebra, where $`G^{}`$ is a Lie group. Indeed, let $`J_a`$ denote the currents of the affine $`G^{}`$ current algebra with central extension $`k^{}`$, and define the generators
$$I_{a;n}=\frac{dz}{2\pi i}\gamma ^n(z)J_a(z)$$
(9)
From the defining OPE
$$J_a(z)J_b(w)=\frac{\kappa _{a,b}k^{}}{(zw)^2}+\frac{f_{a,b}^{}{}_{}{}^{c}J_c(w)}{zw}$$
(10)
where $`\kappa _{a,b}`$ and $`f_{a,b}^{}{}_{}{}^{c}`$ are the Cartan-Killing form and the structure constants, respectively, of the underlying Lie algebra $`𝐠^{}`$, one finds
$`[L_n,I_{b;m}]`$ $`=`$ $`mI_{b;n+m}`$
$`[I_{a;n},I_{b;m}]`$ $`=`$ $`f_{a,b}^{}{}_{}{}^{c}I_{c;n+m}+npk^{}\kappa _{a,b}\delta _{n+m,0}`$ (11)
Summation over “properly” repeated indices is implicit. The central extension $`k^{}p`$ of the space-time affine Lie algebra generated by $`\left\{I_{a;n}\right\}`$ is seen to be the one, $`k^{}`$, of the world sheet $`G^{}`$ current algebra multiplied by the winding number $`p`$ of the embedded $`sl(2)`$ subalgebra. As the currents $`J_a`$ are spin one primary fields, the construction is readily seen to be BRST invariant due to the decomposition $`SL(2)G^{}`$.
## 3 Affine Current Superalgebra
### 3.1 Lie Superalgebra
Let g<sup>0</sup> and g<sup>1</sup> denote the even and odd parts, respectively, of the Lie superalgebra g of rank $`r`$, see Ref. and references therein. $`\mathrm{\Delta }=\mathrm{\Delta }^0\mathrm{\Delta }^1`$ is the set of roots $`\alpha `$ of g where $`\mathrm{\Delta }^0`$ ($`\mathrm{\Delta }^1`$) is the set of even (odd) roots. The set of positive roots $`\alpha >0`$ is $`\mathrm{\Delta }_+=\mathrm{\Delta }_+^0\mathrm{\Delta }_+^1`$. A choice of simple roots is written $`\left\{\alpha _i\right\}_{i=1,\mathrm{},r}`$. A distinguished representation is characterized by exactly one simple root being odd. Related to the triangular decomposition
$$𝐠=𝐠_{}𝐡𝐠_+$$
(12)
the raising and lowering operators are denoted $`E_\alpha ,J_\alpha 𝐠_+`$ and $`F_\alpha ,J_\alpha 𝐠_{}`$, respectively, with $`\alpha \mathrm{\Delta }_+`$, while $`H_i,J_i𝐡`$ are the Cartan generators. Generic Lie superalgebra elements are denoted $`J_a`$ and satisfy
$$[J_a,J_b\}=f_{a,b}^{}{}_{}{}^{c}J_c$$
(13)
where $`[,\}`$ is an anti-commutator if both arguments are fermionic, and otherwise a commutator. The Jacobi identities read
$$[J_a,[J_b,J_c\}\}=[[J_a,J_b\},J_c\}+(1)^{p(J_a)p(J_b)}[J_b,[J_a,J_c\}\}$$
(14)
where the parity $`p(J_a)`$ is 1 (0) for $`J_a`$ an odd (even) generator. The Cartan-Killing form $`\kappa _{a,b}`$
$$2h^{}\kappa _{a,b}=\text{str}(\text{ad}_{J_a}\text{ad}_{J_b})$$
(15)
and the Cartan matrix $`A_{ij}=\alpha _j(H_i)`$ are related as $`\kappa _{i,j}=A_{ij}\kappa _{\alpha _j,\alpha _j}`$. $`h^{}`$ is the dual Coxeter number. The Weyl vector
$`\rho =\rho ^0\rho ^1`$
$`\rho ^0={\displaystyle \frac{1}{2}}{\displaystyle \underset{\alpha \mathrm{\Delta }_+^0}{}}\alpha `$ , $`\rho ^1={\displaystyle \frac{1}{2}}{\displaystyle \underset{\alpha \mathrm{\Delta }_+^1}{}}\alpha `$ (16)
satisfies $`\rho \alpha _i=\alpha _i^2/2`$.
For each positive even or odd root $`\alpha >0`$ we introduce a super-triangular coordinate denoted by $`x^\alpha `$ or $`\theta ^\alpha `$, respectively, where $`\theta ^\alpha `$ is Grassmann odd. In terms of the matrix
$$C_a^b(x,\theta )=\underset{\alpha \mathrm{\Delta }_+^0}{}x^\alpha f_{\alpha ,a}^{}{}_{}{}^{b}\underset{\alpha \mathrm{\Delta }_+^1}{}\theta ^\alpha f_{\alpha ,a}^{}{}_{}{}^{b}$$
(17)
one may then realize the Lie superalgebra in terms of differential operators
$$J_a(x,\theta ,,\mathrm{\Lambda })=\underset{\alpha >0}{}V_a^\alpha (x,\theta )_\alpha +\underset{j=1}{\overset{r}{}}P_a^j(x,\theta )\mathrm{\Lambda }_j$$
(18)
where $`\mathrm{\Lambda }`$ is the weight of the representation, and $`\mathrm{\Lambda }_j`$ are the labels defined by
$$H_j|\mathrm{\Lambda }=\mathrm{\Lambda }(H_j)|\mathrm{\Lambda }=\mathrm{\Lambda }_j|\mathrm{\Lambda }$$
(19)
$`_\alpha `$ is differentiation with respect to $`x^\alpha `$ or $`\theta ^\alpha `$ depending on the parity of $`\alpha `$, whereas $`V`$ and $`P`$ are finite dimensional polynomials:
$`V_\alpha ^\alpha ^{}(x,\theta )`$ $`=`$ $`\left[B(C(x,\theta ))\right]_\alpha ^\alpha ^{}`$
$`V_i^\alpha ^{}(x,\theta )`$ $`=`$ $`\left[C(x,\theta )\right]_i^\alpha ^{}`$
$`V_\alpha ^\alpha ^{}(x,\theta )`$ $`=`$ $`{\displaystyle \underset{\alpha ^{\prime \prime }>0}{}}\left[e^{C(x,\theta )}\right]_\alpha ^{\alpha ^{\prime \prime }}\left[B(C(x,\theta ))\right]_{\alpha ^{\prime \prime }}^\alpha ^{}`$
$`P_\alpha ^j(x,\theta )`$ $`=`$ $`0`$
$`P_i^j(x,\theta )`$ $`=`$ $`\delta _i^j`$
$`P_\alpha ^j(x,\theta )`$ $`=`$ $`\left[e^{C(x,\theta )}\right]_\alpha ^j`$ (20)
$`B(u)`$ is the generating function for the Bernoulli numbers $`B_n`$
$`B(u)`$ $`=`$ $`{\displaystyle \frac{u}{e^u1}}={\displaystyle \underset{n0}{}}{\displaystyle \frac{B_n}{n!}}u^n`$
$`(B(u))^1`$ $`=`$ $`{\displaystyle \frac{e^u1}{u}}={\displaystyle \underset{n0}{}}{\displaystyle \frac{1}{(n+1)!}}u^n`$ (21)
The formal power series expansions (20) all truncate and become polynomials due to the nilpotency of the matrix $`C`$ (17). For later use, let us also introduce the notation $`V_{+}^{}{}_{}{}^{+}`$ for the first of the polynomials in (20)
$$V_{+}^{}{}_{}{}^{+}(x,\theta )=[B(C(x,\theta ))]_{+}^{}{}_{}{}^{+}=B(C_{+}^{}{}_{}{}^{+}(x,\theta ))$$
(22)
where $`C_{+}^{}{}_{}{}^{+}`$ is the submatrix of $`C`$ (17) with both row and column indices positive (even or odd) roots. $`V_{+}^{}{}_{}{}^{+}`$ is immediately seen to be invertible
$$(V_{+}^{}{}_{}{}^{+}(x,\theta ))^1=(B(C_{+}^{}{}_{}{}^{+}(x,\theta )))^1=\underset{n0}{}\frac{1}{(n+1)!}(C_{+}^{}{}_{}{}^{+}(x,\theta ))^n$$
(23)
Most Lie superalgebras with even subalgebra $`𝐠^0=sl(2)𝐠^{}`$ have the property that the embedding of $`sl(2)`$ in $`𝐠`$ carried by $`𝐠^1`$ is a spin 1/2 representation<sup>3</sup><sup>3</sup>3This is true for all basic Lie superalgebras with even subalgebra $`𝐠^0=sl(2)𝐠^{}`$ except $`osp(3|2M)`$ where the embedding is a spin 1 representation, see e.g. .. This means that the space of odd roots may be divided into two parts
$$\mathrm{\Delta }^1=\mathrm{\Delta }^1\mathrm{\Delta }^{1+}$$
(24)
where the roots $`\alpha ^\pm \mathrm{\Delta }^{1\pm }`$ are characterized by
$$\frac{\alpha _{sl(2)}\alpha ^\pm }{\alpha _{sl(2)}^2}=\pm \frac{1}{2}$$
(25)
and we have the correspondence
$$\mathrm{\Delta }^{1+}=\alpha _{sl(2)}+\mathrm{\Delta }^1$$
(26)
$`\alpha _{sl(2)}`$ is the positive root associated to the embedded $`sl(2)`$. In particular, the division (24) is present in the case of our main interest, namely the Lie superalgebra $`sl(2|M)`$ which is considered in Section 4 and further in Appendix A.
### 3.2 Free Field Realization
Associated to a Lie superalgebra is an affine Lie superalgebra characterized by the central extension $`k`$, and associated to an affine Lie superalgebra is an affine current superalgebra whose generators are conformal spin one primary fields and have the mutual operator product expansions<sup>4</sup><sup>4</sup>4We note that the extension of the Virasoro algebra to include an affine Lie algebra (11) discussed in Section 2 may readily be generalized to an affine Lie superalgebra simply by substituting the Lie group $`G^{}`$ with a Lie supergroup; the only change being that the commutator $`[I_{a;n},I_{b;m}]`$ becomes an anti-commutator for $`J_a`$ and $`J_b`$ both fermionic.
$$J_a(z)J_b(w)=\frac{\kappa _{a,b}k}{(zw)^2}+\frac{f_{a,b}^{}{}_{}{}^{c}J_c(w)}{zw}$$
(27)
We use the same notation $`J,E,F,H`$ for the currents as for the algebra generators. Hopefully, this will not lead to misunderstandings. The associated Sugawara construction
$$T=\frac{1}{2(k+h^{})}\kappa ^{a,b}J_aJ_b$$
(28)
generates the Virasoro algebra with central charge
$$c=\frac{k\text{sdim}(\text{g})}{k+h^{}}$$
(29)
The standard free field construction consists in introducing for every positive even root $`\alpha \mathrm{\Delta }_+^0`$, a pair of free bosonic ghost fields ($`\beta _\alpha ,\gamma ^\alpha `$) of conformal weights (1,0) satisfying the OPE
$$\beta _\alpha (z)\gamma ^\alpha ^{}(w)=\frac{\delta _{\alpha }^{}{}_{}{}^{\alpha ^{}}}{zw}$$
(30)
The corresponding energy-momentum tensor is
$$T_{\beta \gamma }=\underset{\alpha \mathrm{\Delta }_+^0}{}\gamma ^\alpha \beta _\alpha $$
(31)
with central charge
$$c_{\beta \gamma }=2|\mathrm{\Delta }_+^0|=\text{dim}(\text{g}^0)r$$
(32)
For every positive odd root $`\alpha \mathrm{\Delta }_+^1`$ one introduces a pair of free fermionic ghost fields ($`b_\alpha ,c^\alpha `$) of conformal weights (1,0) satisfying the OPE
$$b_\alpha (z)c^\alpha ^{}(w)=\frac{\delta _{\alpha }^{}{}_{}{}^{\alpha ^{}}}{zw}$$
(33)
The corresponding energy-momentum tensor is
$$T_{bc}=\underset{\alpha \mathrm{\Delta }_+^1}{}c^\alpha b_\alpha $$
(34)
with central charge
$$c_{bc}=2|\mathrm{\Delta }_+^1|=\text{dim}(\text{g}^1)$$
(35)
For every Cartan index $`i=1,\mathrm{},r`$ one introduces a free scalar boson $`\phi _i`$ with contraction
$$\phi _i(z)\phi _j(w)=\kappa _{i,j}\mathrm{ln}(zw)$$
(36)
The corresponding energy-momentum tensor
$$T_\phi =\frac{1}{2}\phi \phi \frac{1}{\sqrt{k+h^{}}}\rho ^2\phi $$
(37)
has central charge
$$c_\phi =r\frac{h^{}\text{sdim}(\text{g})}{k+h^{}}$$
(38)
where the super-dimension sdim$`(\text{g})`$ of the Lie superalgebra g is defined as the difference dim(g<sup>0</sup>) $``$ dim(g<sup>1</sup>). In obtaining (38) we have used Freudenthal-de Vries (super-)strange formula
$$\rho ^2=\frac{h^{}}{12}\text{sdim}(\text{g})$$
(39)
The total free field realization of the Sugawara energy-momentum tensor is $`T=T_{\beta \gamma }+T_{bc}+T_\phi `$ and has indeed central charge (29).
The generalized Wakimoto free field realization of the affine current superalgebra is obtained by the substitution
$`_{x^\alpha }\beta _\alpha (z)`$ $`,`$ $`x^\alpha \gamma ^\alpha (z),\mathrm{\Lambda }_i\sqrt{k+h^{}}\phi _i(z)`$
$`_{\theta ^\alpha }b_\alpha (z)`$ $`,`$ $`\theta ^\alpha c^\alpha (z)`$ (40)
in the differential operator realization $`\left\{J_a(x,\theta ,,\mathrm{\Lambda })\right\}`$ (18), (20), and a subsequent addition of anomalous terms linear in $`\gamma `$ or $`c`$:
$`J_a(z)`$ $`=`$ $`{\displaystyle \underset{\alpha \mathrm{\Delta }_+^0}{}}V_a^\alpha (\gamma (z),c(z))\beta _\alpha (z)+{\displaystyle \underset{\alpha \mathrm{\Delta }_+^1}{}}V_a^\alpha (\gamma (z),c(z))b_\alpha (z)`$ (41)
$`+`$ $`\sqrt{k+h^{}}{\displaystyle \underset{j=1}{\overset{r}{}}}P_a^j(\gamma (z),c(z))\phi _j(z)+J_a^{\text{anom}}(\gamma (z),c(z),\gamma (z),c(z))`$
Anomalous terms are only added to the lowering generators $`F_\alpha (z)`$
$$J_a^{\text{anom}}(\gamma (z),c(z),\gamma (z),c(z))=\{\begin{array}{c}0\text{for}a=i,\alpha >0\hfill \\ \\ _{\alpha ^{}\mathrm{\Delta }_+^0}\gamma ^\alpha ^{}(z)F_{\alpha ,\alpha ^{}}(\gamma (z),c(z))\hfill \\ +_{\alpha ^{}\mathrm{\Delta }_+^1}c^\alpha ^{}(z)F_{\alpha ,\alpha ^{}}(\gamma (z),c(z))\text{for}a=\alpha <0\hfill \end{array}$$
(42)
and are given by
$`F_{\alpha \mathrm{\Delta }_+^0,\alpha ^{}}(\gamma ,c)`$ $`=`$ $`k{\displaystyle \underset{\mu \mathrm{\Delta }_+}{}}\left[(V_{+}^{}{}_{}{}^{+}(\gamma ,c))^1\right]_\alpha ^{}^\mu \kappa _{\mu ,\alpha }`$
$`+`$ $`{\displaystyle \underset{\mu ,\sigma \mathrm{\Delta }_+^0,\nu \mathrm{\Delta }_+}{}}\left[(V_{+}^{}{}_{}{}^{+}(\gamma ,c))^1\right]_\alpha ^{}^\mu _\sigma V_\mu ^\nu (\gamma ,c)_\nu V_\alpha ^\sigma (\gamma ,c)`$
$``$ $`{\displaystyle \underset{\mu \mathrm{\Delta }_+^0,\sigma \mathrm{\Delta }_+^1,\nu \mathrm{\Delta }_+}{}}\left[(V_{+}^{}{}_{}{}^{+}(\gamma ,c))^1\right]_\alpha ^{}^\mu _\sigma V_\mu ^\nu (\gamma ,c)_\nu V_\alpha ^\sigma (\gamma ,c)`$
$`+`$ $`{\displaystyle \underset{\mu \mathrm{\Delta }_+^1,\sigma ,\nu \mathrm{\Delta }_+}{}}\left[(V_{+}^{}{}_{}{}^{+}(\gamma ,c))^1\right]_\alpha ^{}^\mu _\sigma V_\mu ^\nu (\gamma ,c)_\nu V_\alpha ^\sigma (\gamma ,c)`$
$`F_{\alpha \mathrm{\Delta }_+^1,\alpha ^{}}(\gamma ,c)`$ $`=`$ $`k{\displaystyle \underset{\mu \mathrm{\Delta }_+}{}}\left[(V_{+}^{}{}_{}{}^{+}(\gamma ,c))^1\right]_\alpha ^{}^\mu \kappa _{\mu ,\alpha }`$ (43)
$`+`$ $`{\displaystyle \underset{\mu ,\mathrm{\Delta }_+^0,\sigma ,\nu \mathrm{\Delta }_+}{}}\left[(V_{+}^{}{}_{}{}^{+}(\gamma ,c))^1\right]_\alpha ^{}^\mu _\sigma V_\mu ^\nu (\gamma ,c)_\nu V_\alpha ^\sigma (\gamma ,c)`$
$`+`$ $`{\displaystyle \underset{\mu \mathrm{\Delta }_+^1,\sigma \mathrm{\Delta }_+^0,\nu \mathrm{\Delta }_+}{}}\left[(V_{+}^{}{}_{}{}^{+}(\gamma ,c))^1\right]_\alpha ^{}^\mu _\sigma V_\mu ^\nu (\gamma ,c)_\nu V_\alpha ^\sigma (\gamma ,c)`$
$``$ $`{\displaystyle \underset{\mu ,\sigma \mathrm{\Delta }_+^1,\nu \mathrm{\Delta }_+}{}}\left[(V_{+}^{}{}_{}{}^{+}(\gamma ,c))^1\right]_\alpha ^{}^\mu _\sigma V_\mu ^\nu (\gamma ,c)_\nu V_\alpha ^\sigma (\gamma ,c)`$
In particular, for $`\alpha `$ a simple root the anomalous term $`F_{\alpha ,\alpha ^{}}`$ is a constant independent of $`\gamma `$ and $`c`$:
$$F_{\alpha _i,\alpha ^{}}(\gamma ,c)=\frac{1}{2}\delta _{\alpha _i,\alpha ^{}}\left((2k+h^{})\kappa _{\alpha _i,\alpha _i}A_{ii}\right)$$
(44)
This concludes the explicit free field realization of general affine current superalgebras obtained in Ref. , where the polynomials $`V`$, $`P`$ and $`(V_{+}^{}{}_{}{}^{+})^1`$ are given in (20) and (23).
## 4 Generators of the Superconformal Algebra
In this section we shall construct the Virasoro generators and the supercurrents of the SCA in space-time that may be induced by the affine $`SL(2|N/2)`$ current superalgebra on the world sheet. Some steps towards constructing the complete set of SCA generators are also taken. For simplicity, we consider the underlying Lie superalgebra $`sl(2|N/2)`$ in the distinguished representation, see Appendix A. Let us introduce the abbreviations $`\gamma =\gamma ^{\alpha _1}`$, $`E_1=E_{\alpha _1}`$ etc for the objects related to the embedded $`sl(2)`$. Hopefully, no misunderstandings will arise, as we are also using $`\gamma `$ to represent a general bosonic ghost field argument in the polynomials $`V`$ and $`P`$, though in general we will leave out the arguments.
Using the explicit polynomials
$`V_{\alpha _1}^{\alpha _1}(\gamma ,c)`$ $`=`$ $`1`$
$`V_1^{\alpha _1}(\gamma ,c)`$ $`=`$ $`2\gamma `$
$`V_{\alpha _1}^{\alpha _1}(\gamma ,c)`$ $`=`$ $`\gamma ^2`$ (45)
it is straightforward to verify that the Virasoro algebra (3) is generated by
$`L_n`$ $`=`$ $`{\displaystyle \frac{dz}{2\pi i}_n(z)}`$
$`_n`$ $`=`$ $`a_+(n)\gamma ^{n+1}E_1+a_3(n)\gamma ^nH_1+a_{}(n)\gamma ^{n1}F_1`$ (46)
and has central charge
$$c=6k_1^{}p_1$$
(47)
$`p_1`$ is the winding number (7) for the ghost field $`\gamma ^{\alpha _1}=\gamma `$, and $`k_1^{}=\kappa _{\alpha _1,\alpha _1}k`$ is the level of the embedded $`sl(2)`$ or the level in the direction $`\alpha _1`$.
As a preparation for constructing the algebra generators, let us introduce the generators
$$J_{a;n}=\frac{dz}{2\pi i}\gamma ^n(z)J_a(z)$$
(48)
and consider the OPE
$`_n(z)J_a(w)`$ $`=`$ $`{\displaystyle \frac{1}{zw}}\{(a_+(n)\gamma ^{n+1}f_{\alpha _1,a}^{}{}_{}{}^{c}+a_3(n)\gamma ^nf_{1,a}^{}{}_{}{}^{c}+a_{}(n)\gamma ^{n1}f_{\alpha _1,a}^{}{}_{}{}^{c})J_c`$
$`na_3(n)\gamma ^{n2}V_a^{\alpha _1}(\gamma ^2E_1+\gamma H_1F_1)\}`$
$``$ $`{\displaystyle \frac{1}{(zw)^2}}na_3(n)\gamma ^{n2}(z)\left(\gamma ^2(z)V_{\alpha _1}^\nu (z)+\gamma (z)V_1^\nu (z)V_{\alpha _1}^\nu (z)\right)_\nu V_a^{\alpha _1}(w)`$
Here and in the following summations over repeated indices are implicit. Summations over root indices are meant to be over all positive roots if not otherwise indicated. Actually, this restriction is not necessary as the super-triangular coordinates are defined for positive roots only, i.e. $`_\nu `$ exists only for $`\nu >0`$. Now, due to the structure of the root space (see Appendix A) we immediately obtain
$$[L_n,J_{\alpha ;0}]=0\text{for}\alpha \mathrm{\Delta }^0\left\{\pm \alpha _1\right\}$$
(50)
Likewise, it follows that
$`[L_n,J_{\alpha ^{};0}]`$ $`=`$ $`{\displaystyle \frac{1}{2}}(n1)\gamma ^n\left((n+1)J_{\alpha ^{};n}n\gamma J_{\alpha ^+;n+1}\right)`$
$`[L_n,J_{\alpha ^+;0}]`$ $`=`$ $`{\displaystyle \frac{1}{2}}(n+1)\gamma ^{n1}\left(nJ_{\alpha ^{};n1}(n1)\gamma J_{\alpha ^+;n}\right)`$ (51)
The idea is to use the right hand sides in the construction of the supercurrents. To this end let us consider the general setting where a primary field $`\mathrm{\Phi }`$ of weight $`h`$
$$[L_n,\mathrm{\Phi }_{m+\eta }]=\left((h1)nm\eta \right)\mathrm{\Phi }_{n+m+\eta }$$
(52)
may be obtained as a commutator of the form
$$[L_m,B]=b(m;\eta ;h)\mathrm{\Phi }_{m+\eta }$$
(53)
$`\eta `$ is a possible non-integer shift in the modes, whereas $`b`$ is some function. From the Jacobi identities this function satisfies the recursion relation
$$(nm)b(n+m)=((h1)nm\eta )b(m)((h1)mn\eta )b(n)$$
(54)
allowing the simple solution
$$b(m)=b_0+b_1m,(1h)b_0=\eta b_1$$
(55)
This is precisely of the form we have encountered in (51). Since we want to construct supercurrents of weight 3/2 we should choose $`\eta ZZ+\frac{1}{2}`$, and we define the generators
$`G_{\alpha ^{};n+1/2}`$ $`=`$ $`{\displaystyle \frac{dz}{2\pi i}𝒢_{\alpha ^{};n+1/2}(z)}`$
$`𝒢_{\alpha ^{};n+1/2}`$ $`=`$ $`(n+1)\gamma ^nJ_\alpha ^{}n\gamma ^{n+1}J_{\alpha ^+}`$ (56)
Of course, it still remains to verify that the supercurrents $`G`$ are indeed primary:
$$[L_n,G_{\alpha ^{};m+1/2}]=\left(\frac{1}{2}nm\frac{1}{2}\right)G_{\alpha ^{};n+m+1/2}$$
(57)
However, that follows immediately from the simple computation of the OPE $`_n𝒢_{\alpha ^{};m+1/2}`$, and we have thus constructed half of the supercurrents. Note that both commutators in (51) lead to the same supercurrent as we have
$$G_{\alpha ^{};n+1/2}=\{\begin{array}{c}\frac{2}{n1}[L_n,J_{\alpha ^{};0}]\hfill \\ \\ \frac{2}{n+2}[L_{n+1},J_{\alpha ^+;0}]\hfill \end{array}$$
(58)
This means that $`G`$ may be represented as a commutator for all integer modes $`n`$.
Let us now turn to the construction of the supercurrents $`\overline{G}`$. For $`a=\alpha ^\pm \mathrm{\Delta }_{}^{1\pm }`$ it follows from (LABEL:LJ) that a situation like (58) occurs provided
$$V_{\alpha ^+}^{\alpha _1}=\gamma V_\alpha ^{}^{\alpha _1}$$
(59)
This important relation is proven in Appendix A where also some identities involving $`V_\alpha ^{}^{\alpha _1}`$ are derived. We are led to define the supercurrents $`\overline{G}`$ by
$`\overline{G}_{\alpha ^{};n1/2}`$ $`=`$ $`{\displaystyle \frac{dz}{2\pi i}\overline{𝒢}_{\alpha ^{};n1/2}(z)}`$
$`\overline{𝒢}_{\alpha ^{};n1/2}`$ $`=`$ $`(n1)\gamma ^nJ_\alpha ^{}+n\gamma ^{n1}J_{\alpha ^+}`$ (60)
$``$ $`n(n1)\gamma ^{n2}V_\alpha ^{}^{\alpha _1}\left(\gamma ^2E_1+\gamma H_1F_1\right)`$
$`+`$ $`n(n1)\gamma ^{n2}\left(\gamma ^2V_{\alpha _1}^\nu +\gamma V_1^\nu V_{\alpha _1}^\nu \right)_\nu _\sigma V_\alpha ^{}^{\alpha _1}\gamma ^\sigma `$
$`=`$ $`(n1)\gamma ^nJ_\alpha ^{}+n\gamma ^{n1}J_{\alpha ^+}`$
$``$ $`n(n1)\gamma ^{n2}\left(\mathrm{\Gamma }_{\alpha _1}^\nu b_\nu (k1+h^{}/2)\gamma \gamma ^\sigma \mathrm{\Gamma }_{\alpha _1}^\nu _\nu _\sigma \right)V_\alpha ^{}^{\alpha _1}`$
where
$$\mathrm{\Gamma }_{\alpha _1}^\nu =\gamma ^2V_{\alpha _1}^\nu +\gamma V_1^\nu V_{\alpha _1}^\nu $$
(61)
As indicated in (60), one may show that the upper root index $`\nu `$ is always an odd (and positive) root. The analogue to (58) reads
$$\overline{G}_{\alpha ^{};n1/2}=\{\begin{array}{c}\frac{2}{n+1}[L_n,J_{\alpha ^{};0}]\hfill \\ \\ \frac{2}{n2}[L_{n1},J_{\alpha ^+;0}]\hfill \end{array}$$
(62)
As in the case of the supercurrents $`G`$, there is a free overall scaling. However, in order to produce the conventional prefactor of plus one multiplying the Virasoro generator in the anti-commutator $`\{G,\overline{G}\}`$ (see (73)), the relative factor is fixed<sup>5</sup><sup>5</sup>5 A more commonly used convention is a prefactor of plus two. However, we have found it natural to define $`G`$ and $`\overline{G}`$ without introducing any powers of $`\sqrt{2}`$. To comply with the standard convention is straightforward, though..
Before proving that $`\overline{G}`$ is a primary field of weight 3/2
$$[L_n,\overline{G}_{\alpha ^{};m1/2}]=\left(\frac{1}{2}nm+\frac{1}{2}\right)\overline{G}_{\alpha ^{};n+m1/2}$$
(63)
let us observe the following property of the construction. Consider the Jacobi identity
$$[J_{\delta ;0},[L_n,J_{\pm \alpha ^{};0}]]+[L_n,[J_{\pm \alpha ^{};0},J_{\delta ;0}]]=[J_{\pm \alpha ^{};0},[L_n,J_{\delta ;0}]]=0$$
(64)
where $`\delta \mathrm{\Delta }^0\left\{\pm \alpha _1\right\}`$ is any even (positive or negative) root different from $`\pm \alpha _1`$. From the construction of $`G`$ and $`\overline{G}`$ it follows that<sup>6</sup><sup>6</sup>6Subtleties for $`n=\pm 1,\pm 2`$ are immediately resolved by the alternative commutator representations (58) and (62).
$`G_{\beta ^{};n+1/2}`$ $`=`$ $`[J_{\beta ^{}\alpha ^{};0},G_{\alpha ^{};n+1/2}],\text{for}\beta ^{}\alpha ^{}\mathrm{\Delta }^0\left\{\pm \alpha _1\right\}`$
$`\overline{G}_{\beta ^{};n1/2}`$ $`=`$ $`[J_{\alpha ^{}\beta ^{};0},\overline{G}_{\alpha ^{};n1/2}],\text{for}\beta ^{}\alpha ^{}\mathrm{\Delta }^0\left\{\pm \alpha _1\right\}`$ (65)
Here we have used that
$$f_{\alpha ^{},\beta ^{}\alpha ^{}}^{}{}_{}{}^{\beta ^{}}=f_{\alpha ^{},\alpha ^{}\beta ^{}}^{}{}_{}{}^{\beta ^{}}=1,\text{for}\beta ^{}\alpha ^{}\mathrm{\Delta }^0\left\{\pm \alpha _1\right\}$$
(66)
Besides providing information on the underlying algebraic structure of our construction, the translational property (65) may be used to reduce considerations for general supercurrents to similar ones for the supercurrents $`G_{\alpha _2;n+1/2}`$ and in particular $`\overline{G}_{\alpha _2;n1/2}`$. $`\alpha _2`$ is the only fermionic simple root, see Appendix A. Thus, as a first application we shall prove that $`\overline{G}`$ is primary. From the Jacobi identities we find
$$[L_n,\overline{G}_{\alpha ^{};m1/2}]=[J_{\alpha _2\alpha ^{};0},[L_n,\overline{G}_{\alpha _2;m1/2}]]$$
(67)
leaving us with the task of proving that $`\overline{G}_{\alpha _2}`$ is primary. To that end we work out
$$V_{\alpha _2}^{\alpha _1}=\mathrm{\Gamma }_{\alpha _1}^{\alpha _2}=\frac{1}{2}\gamma c+C$$
(68)
where $`c`$ and $`C`$ are the fermionic ghost fields associated to the odd roots $`\alpha _2`$ and $`\alpha _1+\alpha _2`$, respectively, and the supercurrent becomes
$`\overline{𝒢}_{\alpha _2;m1/2}`$ $`=`$ $`(m1)\gamma ^mJ_{\alpha _2}+m\gamma ^{m1}J_{\alpha _1\alpha _2}`$ (69)
$``$ $`m(m1)\gamma ^{m2}V_{\alpha _2}^{\alpha _1}\left(\gamma ^2E_1+\gamma H_1F_1{\displaystyle \frac{1}{2}}\gamma \right)`$
Now, one may compute the OPE $`_n\overline{𝒢}_{\alpha _2;m1/2}`$ and reduce the result using
$$V_{\alpha _1}^{\alpha _1+\alpha _2}=\frac{1}{2}c,V_1^{\alpha _1+\alpha _2}=C,V_{\alpha _1}^{\alpha _1+\alpha _2}=\frac{1}{2}\gamma V_{\alpha _2}^{\alpha _1}$$
(70)
to the desired commutator
$$[L_n,\overline{G}_{\alpha _2;m1/2}]=\left(\frac{1}{2}nm+\frac{1}{2}\right)\overline{G}_{\alpha _2;n+m1/2}$$
(71)
This concludes the proof of (63) that $`\overline{G}_\alpha ^{}`$ is primary of weight 3/2.
### 4.1 Affine $`SL(N/2)U(1)`$ Current Subalgebra
In order to derive the entire set of generators of the SCA, one should first consider the anti-commutators $`\{G_\alpha ^{},G_\beta ^{}\}`$, $`\{G_\alpha ^{},\overline{G}_\beta ^{}\}`$ and $`\{\overline{G}_\alpha ^{},\overline{G}_\beta ^{}\}`$. It is readily seen that
$$\{G_{\alpha ^{};n+1/2},G_{\beta ^{};m+1/2}\}=0$$
(72)
whereas a rather cumbersome but essentially straightforward computation reveals that
$$\{G_{\alpha ^{};n+1/2},\overline{G}_{\beta ^{};m1/2}\}=\delta _{\alpha ^{},\beta ^{}}L_{n+m}+(nm+1)K_{\alpha ^{};\beta ^{};n+m}+\frac{1}{6}cn(n+1)\delta _{n+m,0}\delta _{\alpha ^{},\beta ^{}}$$
(73)
where the current $`K`$ is defined by
$`K_{\alpha ^{};\beta ^{};n}`$ $`=`$ $`{\displaystyle \frac{dz}{2\pi i}𝒦_{\alpha ^{};\beta ^{};n}(z)}`$
$`𝒦_{\alpha ^{};\beta ^{};n}`$ $`=`$ $`n\gamma ^{n1}V_\beta ^{}^{\alpha _1}\left(\gamma J_{\alpha ^+}J_\alpha ^{}\right)\gamma ^nf_{\alpha ^{},\beta ^{}}^{}{}_{}{}^{c}J_c`$ (74)
$`+`$ $`{\displaystyle \frac{1}{2}}\delta _{\alpha ^{},\beta ^{}}\gamma ^{n1}\left(n\left(\gamma ^2E_1+\gamma H_1F_1\right)\gamma H_1\right)`$
$`+`$ $`n\gamma ^{n1}\left(\gamma V_{\alpha ^+}^\nu V_\alpha ^{}^\nu \right)_\nu _\sigma V_\beta ^{}^{\alpha _1}\gamma ^\sigma `$
There are several ways of representing $`K_{\alpha ^{};\beta ^{};n}`$ of which the following two turn out to be useful
$$K_{\alpha ^{};\beta ^{};n}=\{\begin{array}{c}\frac{1}{n+1}\left(\{G_{\alpha ^{};n+1/2},\overline{G}_{\beta ^{};1/2}\}\delta _{\alpha ^{},\beta ^{}}L_n\right)\hfill \\ \\ \frac{1}{n1}\left(\{G_{\alpha ^{};1/2},\overline{G}_{\beta ^{};n1/2}\}\delta _{\alpha ^{},\beta ^{}}L_n\right)\hfill \end{array}$$
(75)
In particular, they may be used in a straightforward verification that the current $`K`$ is primary of weight 1:
$$[L_n,K_{\alpha ^{};\beta ^{};m}]=mK_{\alpha ^{};\beta ^{};n+m}$$
(76)
In section 6 we shall provide evidence from considering the classical counterpart, that $`\left\{K_{\alpha ^{};\beta ^{};n}\right\}`$ generate an affine $`SL(N/2)U(1)`$ current subalgebra. The number of generators is accordingly
$$|\mathrm{\Delta }_+^1|^2=(N/2)^2=\text{dim}(sl(N/2))+1$$
(77)
A novel feature of our construction is its asymmetry in the two sets of supercurrents $`\left\{G\right\}`$ and $`\left\{\overline{G}\right\}`$, originating in (72) and
$$\{\overline{G}_{\alpha ^{};n1/2},\overline{G}_{\beta ^{};m1/2}\}0,\text{for}\alpha ^{}\beta ^{},nm,n+m1$$
(78)
A proof at the classical level is presented in Section 6, however it is obvious that a result as (78) at the classical level remains true at the quantum level. The right hand side of (78) involves new fields to be introduced in Section 6.
### 4.2 Underlying Lie Superalgebra
Here we shall express the underlying Lie superalgebra in terms of selected modes of the SCA generators. From the Virasoro generator we have
$$E_1=L_1,H_1=2L_0,F_1=L_1$$
(79)
while the supercurrents allow us to write
$`J_\alpha ^{}=G_{\alpha ^{};1/2},`$ $`J_{\alpha ^+}=G_{\alpha ^{};1/2}`$
$`J_\alpha ^{}=\overline{G}_{\alpha ^{};1/2},`$ $`J_{\alpha ^+}=\overline{G}_{\alpha ^{};1/2}`$ (80)
As we have
$`K_{\alpha ^{};\beta ^{};0}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\delta _{\alpha ^{},\beta ^{}}H_1f_{\alpha ^{},\beta ^{}}^{}{}_{}{}^{c}J_c`$
$`\{J_{ϵ_2\delta _u},J_{(ϵ_2\delta _v)}\}`$ $`=`$ $`\delta _{u,v}\left(2H_2{\displaystyle \underset{i=u^{}}{\overset{u}{}}}H_{u^{}+1}\right)+J_{\delta _v\delta _u}`$ (81)
where $`J_{\delta _v\delta _u}`$ is defined only for $`vu`$ (see Appendix A), we find that the remaining $`(N/2)^2`$ Lie superalgebra generators are given by
$`J_{\delta _v\delta _u}`$ $`=`$ $`K_{ϵ_2\delta _u;(ϵ_2\delta _v);0},\text{for}uv`$
$`H_2`$ $`=`$ $`K_{\alpha _2;\alpha _2;0}L_0`$
$`H_i`$ $`=`$ $`K_{ϵ_2\delta _{i1};(ϵ_2\delta _{i1});0}K_{ϵ_2\delta _{i2};(ϵ_2\delta _{i2});0},\text{for}i=3,\mathrm{},N/2+1`$ (82)
### 4.3 $`N=2`$ Superconformal Algebra
For $`N=2`$ the only positive $`\alpha ^{}`$-root is $`\alpha _2`$ (i.e. $`\mathrm{\Delta }_+^1=\{\alpha _2\}`$) and we have the 4 generators
$`_n`$ $`=`$ $`a_+(n)\gamma ^{n+1}E_1+a_3(n)\gamma ^nH_1+a_{}(n)\gamma ^{n1}F_1`$
$`𝒢_{n+1/2}`$ $`=`$ $`(n+1)\gamma ^nJ_{\alpha _2}n\gamma ^{n+1}J_\theta `$
$`\overline{𝒢}_{n1/2}`$ $`=`$ $`(n1)\gamma ^nJ_{\alpha _2}+n\gamma ^{n1}J_\theta `$
$``$ $`n(n1)\gamma ^{n2}V_{\alpha _2}^{\alpha _1}\left(\gamma ^2E_1+\gamma H_1F_1{\displaystyle \frac{1}{2}}\gamma \right)`$
$`𝒦_n`$ $`=`$ $`n\gamma ^{n1}V_{\alpha _2}^{\alpha _1}\left(\gamma J_\theta J_{\alpha _2}\right){\displaystyle \frac{1}{2}}\gamma ^n(H_1+2H_2)`$ (83)
$`+`$ $`{\displaystyle \frac{1}{2}}n\gamma ^{n1}\left(\gamma ^2E_1+\gamma H_1F_1\gamma \right)`$
where $`\theta =\alpha _1+\alpha _2`$. Note that the contribution $`\frac{1}{2}(k+1/2)n\gamma ^{n1}\gamma `$ to $`𝒦_n`$ vanishes upon integration as $`n\frac{dz}{2\pi i}\gamma ^{n1}\gamma =np\delta _{n,0}=0`$. The $`N=2`$ SCA becomes
$`[L_n,L_m]`$ $`=`$ $`(nm)L_{n+m}+{\displaystyle \frac{c}{12}}(n^3n)\delta _{n+m,0}`$
$`[L_n,A_m]`$ $`=`$ $`((h(A)1)nm)A_{n+m},A\{G,\overline{G},K\}`$
$`\{G_{n+1/2},G_{m+1/2}\}`$ $`=`$ $`\{\overline{G}_{n1/2},\overline{G}_{m1/2}\}=0`$
$`\{G_{n+1/2},\overline{G}_{m1/2}\}`$ $`=`$ $`L_{n+m}+(nm+1)K_{n+m}+{\displaystyle \frac{1}{6}}cn(n+1)\delta _{n+m,0}`$
$`[K_n,G_{m+1/2}]`$ $`=`$ $`{\displaystyle \frac{1}{2}}G_{n+m+1/2},[K_n,\overline{G}_{m1/2}]={\displaystyle \frac{1}{2}}\overline{G}_{n+m1/2}`$
$`[K_n,K_m]`$ $`=`$ $`{\displaystyle \frac{1}{12}}cn\delta _{n+m,0}`$ (84)
and closure is seen to be ensured by the 4 generators (83). This result has already been obtained by Ito , though his construction is based on a slightly different but equivalent free field realization of the associated affine $`SL(2|1)`$ current superalgebra.
## 5 BRST Invariance
Before turning to the classical SCA let us discuss an important property of our construction. As pointed out in Ref. , BRST invariance of the construction of the space-time conformal algebra from a world sheet $`SL(2)`$ current algebra requires the Virasoro generators to be primary fields of weight one with respect to the world sheet energy-momentum tensor. This ensures that the integrated fields commute with the world sheet Virasoro algebra. This requirement carries over to the superconformal case, where all (super-)currents (the Virasoro generators $``$, the supercurrents $`𝒢`$ and $`\overline{𝒢}`$, the affine Lie algebra generators $`𝒦`$ etc) are primary of weight one with respect to the Sugawara energy-momentum tensor of the world sheet $`SL(2|N/2)`$ current superalgebra. A naive inspection immediately tells that the four types of currents considered so far have weight one, so all we need to verify is that they are primary. This amounts to verifying that third and higher order poles in the OPEs with the Sugawara tensor $`T`$ all vanish. From the free field realization of $`T`$ it follows that no higher order poles than third order appears. Using that the affine currents $`J`$ are primary fields, the BRST invariance of the supercurrents $`G`$ is readily confirmed as $`V_{\alpha ^\pm }^{\alpha _1}=0`$. The BRST invariance of the Virasoro generators $`L`$ follows from (45), whereas the invariance of the supercurrents $`\overline{G}`$ is ensured provided
$`0`$ $`=`$ $`n(n1)\gamma ^{n1}V_\alpha ^{}^{\alpha _1}+n(n1)\gamma ^{n2}V_{\alpha ^+}^{\alpha _1}`$ (85)
$``$ $`n(n1)V_\alpha ^{}^{\alpha _1}\left(n\gamma ^{n1}V_{\alpha _1}^{\alpha _1}+(n1)\gamma ^{n2}V_1^{\alpha _1}(n2)\gamma ^{n3}V_{\alpha _1}^{\alpha _1}\right)`$
$``$ $`n(n1)\left(\gamma ^nV_{\alpha _1}^\nu +\gamma ^{n1}V_1^\nu \gamma ^{n2}V_{\alpha _1}^\nu \right)_\nu V_\alpha ^{}^{\alpha _1}`$
The three lines vanish separately due to (59), (45) and (130), respectively. Likewise, BRST invariance of the affine Lie algebra generators $`K`$ amounts to verifying
$`0`$ $`=`$ $`n\gamma ^{n1}f_{\alpha ^{},\beta ^{}}^{}{}_{}{}^{c}V_c^{\alpha _1}`$ (86)
$`+`$ $`{\displaystyle \frac{1}{2}}\delta _{\alpha ^{},\beta ^{}}\left\{n\left((n+1)\gamma ^nV_{\alpha _1}^{\alpha _1}+n\gamma ^{n1}V_1^{\alpha _1}(n1)\gamma ^{n2}V_{\alpha _1}^{\alpha _1}\right)n\gamma ^{n1}V_1^{\alpha _1}\right\}`$
Here we have used that $`V_{\alpha ^\pm }^{\alpha _1}=0`$, and the identity (86) follows from (45) and (132).
One may take a more general point of view observing that the (anti-)commutator of two BRST invariant fields commutes with the world sheet Virasoro generators. This follows from the Jacobi identities. Thus, having established that all but one field appearing on the right hand side of a (anti-)commutator (of two BRST invariant fields) are BRST invariant, is sufficient to conclude that the final field is likewise BRST invariant. A trivial example is the alternative deduction that $`K`$ is BRST invariant following from the anti-commutator (73).
## 6 Classical Superconformal Algebra
Here we shall distinguish between classical and quantum SCAs. Our use of the notion quantum is not in the quantum group sense of $`q`$-deformations but rather as opposed to classical as described in the following. Let us recall the situation for free field realizations of affine current superalgebras discussed in Section 3. In that case one may start with a first order linear differential operator realization of the underlying Lie superalgebra. The free field realization of the associated current superalgebra is then obtained by substituting with (normal ordered products of) free fields (40) and subsequently adding “quantum corrections”, “anomalous terms” or “normal ordering terms” (41). We shall denote the differential operator realization a classical limit or version of the associated “quantum” free field realization. A classical algebra thus defined has vanishing central extensions<sup>7</sup><sup>7</sup>7It should be stressed that a non-vanishing central charge of a classical Virasoro algebra may well exist when classical is defined to denote single contractions only, as $`^2\phi (z)^2\phi (w)=12/(zw)^4`$. Here we have used the convention $`\phi (z)\phi (w)=2\mathrm{ln}(zw)`$. However, terms like $`^2\phi `$ are excluded in our “differential operator realization picture” employed in the present paper. Note that BRST invariance is used as an implicit guideline as $`^2\phi `$ has weight 2 with respect to the world sheet energy-momentum tensor..
A similar situation may be expected in the present case. Thus, there should exist a classical counterpart of the full SCA which allows a differential operator realization (and accordingly has vanishing central extensions). Based on this assumption our program is to first work out the classical SCA for then to perform the appropriate substitutions and additions of anomalous terms in order to obtain the full (quantum) SCA. It should be stressed that to each mode of the generators of the quantum SCA, there is an associated differential operator. This results in an infinite dimensional algebra of differential operators contrary to the situation described above where the classical algebra is a standard (finite dimensional) Lie superalgebra.
Having identified the differential operator
$$A(x,\theta ,,\mathrm{\Lambda })=\underset{a}{}Y^a(x,\theta )J_a(x,\theta ,,\mathrm{\Lambda })$$
(87)
as a generator of the classical SCA, the corresponding quantum generator
$$A=\frac{dz}{2\pi i}𝒜(z)$$
(88)
is obtained by performing the substitutions (40) in $`A(x,\theta ,,\mathrm{\Lambda })`$ and adding appropriate anomalous terms linear in derivatives of the spin 0 ghost fields in order to produce $`𝒜(z)`$:
$`𝒜(z)`$ $`=`$ $`{\displaystyle \underset{a}{}}Y^a(\gamma (z),c(z))J_a(z)`$ (89)
$`+`$ $`{\displaystyle \underset{\alpha \mathrm{\Delta }_+^0}{}}X_\alpha (\gamma (z),c(z))\gamma ^\alpha (z)+{\displaystyle \underset{\alpha \mathrm{\Delta }_+^1}{}}X_\alpha ^{}(\gamma (z),c(z))c^\alpha (z)`$
So with the ansatz (87), $`𝒜(z)`$ is linear in the affine currents $`J_a(z)`$ with spin 0 ghost field dependent coefficients. Note that in the expression (89) some anomalous terms are “hidden” in the definition of $`J_a(z)`$, cf. (41).
The question of BRST invariance of $`A`$ (88) may be addressed even without explicit knowledge on the anomalous term. This follows from the fact that any term of the form $`_\alpha Z_\alpha (\gamma ,c)\gamma ^\alpha +_\alpha Z_\alpha ^{}(\gamma ,c)c^\alpha `$ is primary of weight one. Thus, in order to establish that $`A`$ is BRST invariant it suffices to consider the term linear in the affine (super-)currents. In the following we shall accordingly define a classical differential operator to be BRST invariant when its “naively quantized” form linear in the affine (super-)currents is BRST invariant.
Before continuing our program let us briefly justify it. It has turned out to be an immense technical task to complete the derivation of the SCA for general $`N`$. Even at the classical level, the computations are rather involved. A study of the center-less classical SCA seems therefore a natural first project to concentrate on, and one from which one may get structural insight into the full quantum SCA. In the following we shall present some essential steps in the direction of deriving the classical SCA. In the presumably most interesting case of $`N=4`$, the classical SCA is completed. The full quantum $`N=4`$ with generic central charge will be presented elsewhere .
### 6.1 Algebra Generators
In the remaining part of this Section all fields $`A`$ are represented by their classical differential operator analogues $`A(x,\theta ,,\mathrm{\Lambda })`$. To be explicit, let us summarize our findings for the classical generators:
$`L_n`$ $`=`$ $`a_+(n)x^{n+1}E_1+a_3(n)x^nH_1+a_{}(n)x^{n1}F_1`$
$`G_{\alpha ^{};n+1/2}`$ $`=`$ $`(n+1)x^nJ_\alpha ^{}nx^{n+1}J_{\alpha ^+}`$
$`\overline{G}_{\alpha ^{};n1/2}`$ $`=`$ $`(n1)x^nJ_\alpha ^{}+nx^{n1}J_{\alpha ^+}`$
$``$ $`n(n1)x^{n2}V_\alpha ^{}^{\alpha _1}\left(x^2E_1+xH_1F_1\right)`$
$`K_{\alpha ^{};\beta ^{};n}`$ $`=`$ $`nx^{n1}V_\beta ^{}^{\alpha _1}\left(xJ_{\alpha ^+}J_\alpha ^{}\right)x^nf_{\alpha ^{},\beta ^{}}^{}{}_{}{}^{c}J_c`$
$`+`$ $`{\displaystyle \frac{1}{2}}\delta _{\alpha ^{},\beta ^{}}x^{n1}\left(nx^2E_1+(n1)xH_1nF_1\right)`$
$`h(L,G,\overline{G},K)`$ $`=`$ $`(2,3/2,3/2,1)`$ (90)
$`J_a`$ ($`E`$, $`H`$ or $`F`$) denotes the differential operator $`J_a(x,\theta ,,\mathrm{\Lambda })`$ given in (18), (20) while $`V_\beta ^{}^{\alpha _1}`$ is the polynomial in the super-triangular coordinates $`x`$ and $`\theta `$ given in (20). Here and in the following $`x`$ may denote either $`x^{\alpha _1}`$ or a general triangular coordinate argument, though it should be clear from the context which it is. $`h(A_m)`$ indicates that $`A_m`$ is primary of weight $`h`$:
$$[L_n,A_m]=((h(A)1)nm)A_{n+m}$$
(91)
The generators respect among others the anti-commutators
$`\{G_{\alpha ^{};n+1/2},G_{\beta ^{};m+1/2}\}`$ $`=`$ $`0`$
$`\{G_{\alpha ^{};n+1/2},\overline{G}_{\beta ^{};m1/2}\}`$ $`=`$ $`\delta _{\alpha ^{},\beta ^{}}L_{n+m}+(nm+1)K_{\alpha ^{};\beta ^{};n+m}`$ (92)
We shall now discuss the subalgebra generated by $`\left\{K\right\}`$. One finds straightforwardly
$$[K_{\alpha ^{};\beta ^{};n},K_{\mu ^{};\nu ^{};m}]=\delta _{\mu ^{},\beta ^{}}K_{\alpha ^{};\nu ^{};n+m}\delta _{\alpha ^{},\nu ^{}}K_{\mu ^{};\beta ^{};n+m}$$
(93)
In order to show explicitly that this has the affine structure
$$SL(N/2)U(1)$$
(94)
we introduce the following notation. From Appendix A we know that any root $`\alpha ^{}\mathrm{\Delta }_+^1`$ may be represented as $`ϵ_2\delta _u`$ for some $`u=1,\mathrm{},N/2`$, so abbreviate $`K`$ by
$$K_{u;v;n}=K_{ϵ_2\delta _u;(ϵ_2\delta _v);n}$$
(95)
Note also that $`\delta _v\delta _u>0`$ for $`u>v`$. Define now
$`\stackrel{~}{E}_{i;n}`$ $`=`$ $`K_{i+1;i;n}`$
$`\stackrel{~}{H}_{i;n}`$ $`=`$ $`K_{i+1;i+1;n}K_{i;i;n}`$
$`\stackrel{~}{F}_{i;n}`$ $`=`$ $`K_{i;i+1;n}`$ (96)
where $`i=1,\mathrm{},N/21`$ by construction. One may then show that these correspond to the Chevalley generators of an (center-less) affine $`SL(N/2)`$ Lie algebra. In general, the currents $`K_{u;v;n}`$ correspond to raising operators for $`u>v`$, and to lowering operators for $`u<v`$. Furthermore, the generator
$$U_n=\underset{u=1}{\overset{N/2}{}}K_{u;u;n}$$
(97)
is seen to commute with all ladder operators $`K_{u;vu;m}`$, with the Cartan generators $`\stackrel{~}{H}_{i,m}`$ and with $`U_m`$ itself. Thus, $`U`$ generates a (center-less) $`U(1)`$ current algebra and we have the decomposition (94).
Let us return to the anti-commutator $`\{\overline{G}_\alpha ^{},\overline{G}_\beta ^{}\}`$ and prove the classical counterpart of the assertion (78). The anti-commutator may be computed directly or obtained as a special case of a much more general consideration: Introduce the operator
$`\overline{G}_{\beta _1^{},\mathrm{},\beta _k^{};m1+k/2}`$ $`=`$ $`V_{\beta _1^{}}^{\alpha _1}\mathrm{}V_{\beta _{k1}^{}}^{\alpha _1}\left\{(m2+k)x^mJ_{\beta _k^{}}+mx^{m1}J_{\beta _k^+}\right\}`$ (98)
$`+`$ $`\mathrm{}`$
$`\mathrm{}`$
$`+`$ $`\left\{(m2+k)x^mJ_{\beta _1^{}}+mx^{m1}J_{\beta _1^+}\right\}V_{\beta _2^{}}^{\alpha _1}\mathrm{}V_{\beta _k^{}}^{\alpha _1}`$
$``$ $`V_{\beta _1^{}}^{\alpha _1}\mathrm{}V_{\beta _k^{}}^{\alpha _1}\{(m+k1)(m+k2)x^mE_1`$
$`+(m+k2)mx^{m1}H_1m(m1)x^{m2}F_1\}`$
which is seen to reduce to $`\overline{G}_\beta ^{}`$ for $`k=1`$. $`\overline{G}_{\beta _1^{},\mathrm{},\beta _k^{};m}`$ is bosonic (fermionic) for $`k`$ even (odd). In the latter notation the mode $`m`$ is meant to be integer or half-integer depending on the parity of the generator, i.e. depending on $`k`$ being even or odd, respectively. Note that $`\overline{G}_{\beta _1^{},\mathrm{},\beta _k^{}}`$ is anti-symmetric in its root indices. $`J_{\beta _j^\pm }`$ is defined not to act on $`V_{\beta _{j+1}^{}}^{\alpha _1}\mathrm{}V_{\beta _k^{}}^{\alpha _1}`$, and is only written to the left of the $`V`$-monomial for convenience of notation. Thus, within $`\overline{G}_{\beta _1^{},\mathrm{},\beta _k^{}}`$ one has $`J_{\beta _j^\pm }V_{\beta _{j+1}^{}}^{\alpha _1}\mathrm{}V_{\beta _k^{}}^{\alpha _1}=(1)^{kj}V_{\beta _{j+1}^{}}^{\alpha _1}\mathrm{}V_{\beta _k^{}}^{\alpha _1}J_{\beta _j^\pm }`$. This resembles normal ordering needed in the free field realization and may be regarded as a normal ordering of the differential operator. It is to be employed throughout this section. One may now work out the (anti-)commutator
$$[\overline{G}_{\beta _1^{},\mathrm{},\beta _k^{};n},\overline{G}_{\lambda _1^{},\mathrm{},\lambda _l^{};m}\}=\left((k2)m(l2)n\right)\overline{G}_{\beta _1^{},\mathrm{},\beta _k^{},\lambda _1^{},\mathrm{},\lambda _l^{};n+m}$$
(99)
We observe that $`\left\{\overline{G}_{\beta _1^{},\mathrm{},\beta _k^{}}\right\}_{k=1}^{|\mathrm{\Delta }_{}^1|}`$ generate a subalgebra of dimension $`2^{|\mathrm{\Delta }_{}^1|}1=2^{N/2}1`$ and that $`2^{N/21}`$ of the generators are fermionic. Note that the commutator for $`k=l=2`$ vanishes identically. One may also show that (classically) the generators are primary:
$$[L_n,\overline{G}_{\beta _1^{},\mathrm{},\beta _k^{};m}]=((1k/2)nm)\overline{G}_{\beta _1^{},\mathrm{},\beta _k^{};n+m},h=2k/2$$
(100)
Thus, including $`L`$ as a generator of the subalgebra it has dimension $`2^{|\mathrm{\Delta }_{}^1|}`$ and equal numbers of bosonic and fermionic generators. BRST invariance is readily verified, either directly or as a consequence of the recursive relation (99) and the general approach of Section 5.
We note that the generator (98) may be written in the following compact form
$`\overline{G}_{\beta _1^{},\mathrm{},\beta _k^{};m1+k/2}`$ $`=`$ $`x^{m2}\{{\displaystyle \underset{j=1}{\overset{k}{}}}((m+k2)x^2J_{\beta _j^{}}+mxJ_{\beta _j^+}){\displaystyle \frac{}{V_{\beta _j^{}}^{\alpha _1}}}`$ (101)
$`(m+k1)(m+k2)x^2E_1(m+k2)mxH_1`$
$`+m(m1)F_1\}(V_{\beta _1^{}}^{\alpha _1}\mathrm{}V_{\beta _k^{}}^{\alpha _1})`$
where we have defined
$$\frac{}{V_{\beta _j^{}}^{\alpha _1}}V_{\beta _i^{}}^{\alpha _1}\delta _{ij}$$
(102)
As $`\overline{G}_{\beta _1^{},\mathrm{},\beta _k^{}}`$ is a first order differential operator, $`/V_{\beta _j^{}}^{\alpha _1}`$ is meant to act only on the explicitly written products of $`V`$’s within $`\overline{G}_{\beta _1^{},\mathrm{},\beta _k^{}}`$ itself. A special situation occurs for $`k=2`$ since we then have
$`\overline{G}_{\beta _1^{},\beta _2^{};m}`$ $`=`$ $`mS_{\beta _1^{},\beta _2^{};m}`$
$`S_{\beta _1^{},\beta _2^{};m}`$ $`=`$ $`x^{m2}\{{\displaystyle \underset{j=1}{\overset{2}{}}}(x^2J_{\beta _j^{}}+xJ_{\beta _j^+}){\displaystyle \frac{}{V_{\beta _j^{}}^{\alpha _1}}}`$ (103)
$`(m+1)x^2E_1mxH_1+(m1)F_1\}(V_{\beta _1^{}}^{\alpha _1}V_{\beta _2^{}}^{\alpha _1})`$
and $`S_{\beta _1^{},\beta _2^{}}`$ is readily seen to have weight 0. $`\overline{G}_{\beta _1^{},\beta _2^{}}`$ may be interpreted as the derivative of the scalar $`S_{\beta _1^{},\beta _2^{}}`$.
The list of generators presented hitherto is by no means exhaustive. Let us consider generators which may be obtained by the adjoint action of $`\left\{G_\alpha ^{}\right\}`$ on $`\overline{G}_{\beta _1^{},\mathrm{},\beta _k^{}}`$. Firstly, we find the (anti-)commutator
$`[G_{\alpha ^{};n+1/2},\overline{G}_{\beta _1^{},\mathrm{},\beta _k^{};m}\}`$ $`=`$ $`(m(k2)n)\mathrm{\Phi }_{\alpha ^{};\beta _1^{},\mathrm{},\beta _k^{};n+m+1/2}`$ (104)
$`+`$ $`{\displaystyle \frac{k2}{k3}}{\displaystyle \underset{j=1}{\overset{k}{}}}(1)^{j1}\delta _{\alpha ^{},\beta _j^{}}\overline{G}_{\beta _1^{},\mathrm{},\widehat{\beta _j^{}},\mathrm{},\beta _k^{};n+m+1/2}`$
where $`\mathrm{\Phi }_{\alpha ^{};\beta _1^{},\mathrm{},\beta _k^{}}`$ corresponds to the special case $`l=1`$ in the following general expression
$`\mathrm{\Phi }_{\mu _1^{},\mathrm{},\mu _l^{};\beta _1^{},\mathrm{},\beta _k^{};m}`$ (105)
$`=`$ $`x^{m(kl)/21}{\displaystyle \underset{i=1}{\overset{l}{}}}(1)^{i1}\{(m(kl)/2)J_{\mu _i^{}}(m+(kl)/2)xJ_{\mu _i^+}`$
$`x{\displaystyle \underset{j=1}{\overset{k}{}}}f_{\mu _i^{},\beta _j^{}}^{}{}_{}{}^{c}J_c{\displaystyle \frac{}{V_{\beta _j^{}}^{\alpha _1}}}\}{\displaystyle \frac{}{V_{\mu _1^{}}^{\alpha _1}}}\mathrm{}\widehat{{\displaystyle \frac{}{V_{\mu _i^{}}^{\alpha _1}}}}\mathrm{}{\displaystyle \frac{}{V_{\mu _l^{}}^{\alpha _1}}}(V_{\beta _1^{}}^{\alpha _1}\mathrm{}V_{\beta _k^{}}^{\alpha _1})`$
$`+`$ $`{\displaystyle \frac{l}{kl2}}x^{m(kl)/21}\{{\displaystyle \underset{j=1}{\overset{k}{}}}(x^2J_{\beta _j^{}}+xJ_{\beta _j^+}){\displaystyle \frac{}{V_{\beta _j^{}}^{\alpha _1}}}`$
$`(m+(kl)/2)x^2E_1(m+(kl)/21)xH_1`$
$`+(m(kl)/2)F_1\}{\displaystyle \frac{}{V_{\mu _1^{}}^{\alpha _1}}}\mathrm{}{\displaystyle \frac{}{V_{\mu _l^{}}^{\alpha _1}}}(V_{\beta _1^{}}^{\alpha _1}\mathrm{}V_{\beta _k^{}}^{\alpha _1})`$
These generators are only defined for certain integer pairs $`(l,k)`$ to be discussed below. A hat over an object indicates that the object is left out. We observe that $`\mathrm{\Phi }_{\mu _1^{},\mathrm{},\mu _l^{};\beta _1^{},\mathrm{},\beta _k^{}}`$ is bosonic (fermionic) for $`l+k`$ even (odd), and that it is anti-symmetric in the positive root indices and in the negative root indices, separately. Using the polynomial relations listed at the end of Appendix A, one may show that $`\mathrm{\Phi }`$ satisfies the recursion relation
$`[G_{\nu ^{};n+1/2},\mathrm{\Phi }_{\mu _1^{},\mathrm{},\mu _l^{};\beta _1^{},\mathrm{},\beta _k^{};m}\}`$ (106)
$`=`$ $`{\displaystyle \frac{k2}{kl2}}{\displaystyle \underset{j=1}{\overset{k}{}}}(1)^{lj+1}\delta _{\nu ^{},\beta _j^{}}\mathrm{\Phi }_{\mu _1^{},\mathrm{},\mu _l^{};\beta _1^{},\mathrm{},\widehat{\beta _j^{}},\mathrm{},\beta _k^{};n+m+1/2}`$
$``$ $`{\displaystyle \frac{l}{kl2}}\mathrm{\Phi }_{\nu ^{},\mu _1^{},\mathrm{},\mu _l^{};\beta _1^{},\mathrm{},\beta _k^{};n+m+1/2}`$
In addition, $`\mathrm{\Phi }_{\mu _1^{},\mathrm{},\mu _l^{};\beta _1^{},\mathrm{},\beta _k^{}}`$ may be shown to be primary of weight
$$h(\mathrm{\Phi }_{\mu _1^{},\mathrm{},\mu _l^{};\beta _1^{},\mathrm{},\beta _k^{}})=1(kl)/2$$
(107)
As for $`\overline{G}_{\beta _1^{},\mathrm{},\beta _k^{}}`$ (98), BRST invariance of $`\mathrm{\Phi }_{\mu _1^{},\mathrm{},\mu _l^{};\beta _1^{},\mathrm{},\beta _k^{}}`$ is verified straightforwardly, either directly or indirectly. There is a slight subtlety in (106) for $`k=l=1`$, which is easily resolved, though, as $`\mathrm{\Phi }_\mu ^{}`$ (105) may be interpreted as $`G_\mu ^{}`$. One then has
$$[K_{\mu ^{};\beta ^{};n},G_{\nu ^{};m+1/2}]=\delta _{\nu ^{},\beta ^{}}G_{\mu ^{};n+m+1/2}\frac{1}{2}\delta _{\mu ^{},\beta ^{}}G_{\nu ^{};n+m+1/2}$$
(108)
In the definition of $`\mathrm{\Phi }_{\mu _1^{},\mathrm{},\mu _l^{};\beta _1^{},\mathrm{},\beta _k^{}}`$ we obviously have $`kl20`$, and as the expression is obtained by the adjoint action of $`G`$ on $`\overline{G}_{\beta _1^{},\mathrm{},\beta _k^{}}`$, $`1l<k2`$ is seen to be a valid domain. It is relevant for the discussion below on $`N=4`$ that besides $`(l,k)=(1,1)`$ which corresponds to the $`K`$ generators, also $`(l,k)\{(2,1),(1,2),(2,2),(3,2)\}`$ may be reached by this adjoint action. It turns out that three of these four generators may be expressed in terms of simpler generators. First we note that up to permutations in the root indices, $`\mathrm{\Phi }_{\mu ^{},\beta _2^{},\beta _1^{};\beta _1^{},\beta _2^{}}`$ is the only non-vanishing generator of type $`(3,2)`$ (provided $`\mu ^{}\beta _i^{}`$ and $`\beta _1^{}\beta _2^{}`$, of course), while $`\mathrm{\Phi }_{\mu ^{},\beta ^{};\beta ^{}}`$ is the only non-vanishing generator of type $`(2,1)`$ (provided $`\mu ^{}\beta ^{}`$). It is easily shown that they satisfy
$`\mathrm{\Phi }_{\mu _1^{},\mu _2^{},\mu _3^{};\beta _1^{},\beta _2^{};m+1/2}`$ $`=`$ $`\left(\delta _{\mu _3^{},\beta _1^{}}\delta _{\mu _2^{},\beta _2^{}}\delta _{\mu _3^{},\beta _2^{}}\delta _{\mu _2^{},\beta _1^{}}\right)G_{\mu _1^{};m+1/2}`$
$``$ $`\left(\delta _{\mu _3^{},\beta _1^{}}\delta _{\mu _1^{},\beta _2^{}}\delta _{\mu _3^{},\beta _2^{}}\delta _{\mu _1^{},\beta _1^{}}\right)G_{\mu _2^{};m+1/2}`$
$`+`$ $`\left(\delta _{\mu _2^{},\beta _1^{}}\delta _{\mu _1^{},\beta _2^{}}\delta _{\mu _2^{},\beta _2^{}}\delta _{\mu _1^{},\beta _1^{}}\right)G_{\mu _3^{};m+1/2}`$
$`\mathrm{\Phi }_{\mu _1^{},\mu _2^{};\beta ^{};m+1/2}`$ $`=`$ $`\delta _{\mu _2^{},\beta ^{}}G_{\mu _1^{}}\delta _{\mu _1^{},\beta ^{}}G_{\mu _2^{}}`$ (109)
Secondly, we observe that $`\mathrm{\Phi }_{\mu ^{},\alpha ^{};\alpha ^{},\beta ^{}}`$ is the only non-vanishing generator of type $`(2,2)`$ and that it may be reduced as
$`\mathrm{\Phi }_{\mu _1^{},\mu _2^{};\beta _1^{},\beta _2^{};m}`$ $`=`$ $`\delta _{\mu _1^{},\beta _1^{}}K_{\mu _2;\beta _2^{};m}+\delta _{\mu _1^{},\beta _2^{}}K_{\mu _2;\beta _1^{};m}`$ (110)
$`+`$ $`\delta _{\mu _2^{},\beta _1^{}}K_{\mu _1;\beta _2^{};m}\delta _{\mu _2^{},\beta _2^{}}K_{\mu _1;\beta _1^{};m}`$
For $`l>k+1`$, $`k=1,2`$, the generator $`\mathrm{\Phi }_{\mu _1^{},\mathrm{},\mu _l^{};\beta _1^{},\mathrm{},\beta _k^{}}`$ is readily seen to vanish. Relevant for the discussion on the $`N=4`$ SCA in the following, is the result
$`[K_{\mu ^{};\nu ^{};n},\overline{G}_{\alpha ^{};m1/2}]`$ $`=`$ $`{\displaystyle \frac{1}{2}}\delta _{\mu ^{},\nu ^{}}\overline{G}_{\alpha ^{};n+m1/2}\delta _{\mu ^{},\alpha ^{}}\overline{G}_{\nu ^{};n+m1/2}`$ (111)
$`+`$ $`n\mathrm{\Phi }_{\mu ^{};\nu ^{},\alpha ^{};n+m1/2}`$
which one may work out explicitly.
Finally, let us add a comment on the closure of the algebra. Based on the results obtained so far one should expect that the SCA is generated by a BRST invariant set $`\{\stackrel{~}{\mathrm{\Phi }}_{\mu _1^{},\mathrm{},\mu _l^{};\beta _1^{},\mathrm{},\beta _k^{}}\}`$ where $`0l,k`$ and $`lk+1`$. The notation implies for example $`\stackrel{~}{\mathrm{\Phi }}_\mu =G_\mu `$ and $`\stackrel{~}{\mathrm{\Phi }}_{\beta _1^{},\mathrm{},\beta _k^{}}=\overline{G}_{\beta _1^{},\mathrm{},\beta _k^{}}`$, while for $`l=k=0`$ we have $`\stackrel{~}{\mathrm{\Phi }}=L`$. Not all the fields $`\stackrel{~}{\mathrm{\Phi }}`$ need be present as independent fields as illustrated by the reductions (109) and (110). At present we do not have a complete proof of the assumption that $`\{\stackrel{~}{\mathrm{\Phi }}\}`$ generates the SCA, though it is easy to prove (using the polynomial relations of Appendix A) that the following 8 BRST non-invariant “building blocks”
$`V_{\beta _1^{}}^{\alpha _1}\mathrm{}V_{\beta _{k_+^{}}^{}}^{\alpha _1}x^{n_+^{}}J_\alpha ^{},V_{\beta _1^{}}^{\alpha _1}\mathrm{}V_{\beta _{k_+^+}^{}}^{\alpha _1}x^{n_+^+}J_{\alpha ^+},V_{\beta _1^{}}^{\alpha _1}\mathrm{}\widehat{V_{\beta _j^{}}^{\alpha _1}}\mathrm{}V_{\beta _{k_c}^{}}^{\alpha _1}x^{n_c}f_{\alpha ^{},\beta _j^{}}^{}{}_{}{}^{c}J_c`$
$`V_{\beta _1^{}}^{\alpha _1}\mathrm{}V_{\beta _k_{}^{}^{}}^{\alpha _1}x^n_{}^{}J_\alpha ^{},V_{\beta _1^{}}^{\alpha _1}\mathrm{}V_{\beta _{k_{}^+}^{}}^{\alpha _1}x^{n_{}^+}J_{\alpha ^+}`$
$`V_{\beta _1^{}}^{\alpha _1}\mathrm{}V_{\beta _{k_E}^{}}^{\alpha _1}x^{n_E}E_1,V_{\beta _1^{}}^{\alpha _1}\mathrm{}V_{\beta _{k_H}^{}}^{\alpha _1}x^{n_H}H_1,V_{\beta _1^{}}^{\alpha _1}\mathrm{}V_{\beta _{k_F}^{}}^{\alpha _1}x^{n_F}F_1`$ (112)
close under (anti-)commutations.
Despite the fact that closure of the BRST invariant algebra is still not ensured, nor are all the BRST invariant generators at hand, above we have presented substantial evidence that our construction does produce a finitely generated and BRST invariant $`N`$ extended SCA. Below we shall demonstrate that this is indeed the case for $`N=4`$. We intend to come bask elsewhere with further discussion on the BRST invariant SCA and its quantization as described above.
As $`|\mathrm{\Delta }_{}^1|=N/2`$, we observe that for $`N=6,8,\mathrm{}`$ the SCA contains primary generators of negative weight causing problems for the unitarity of the associated CFT, and in particular for its applications to string theory. From an algebraic point of view, however, we see no severe obstacles arising from the appearance of negative weights, and believe that further studies of these algebras are warranted. Their explicit realizations and linearity in the affine currents seem to make such investigations feasible.
### 6.2 Classical $`N=4`$ Superconformal Algebra
We recall that $`sl(2|2)`$ has precisely two positive (fermionic) $`\alpha ^{}`$-roots, see Appendix A:
$$\mathrm{\Delta }_+^1=\{\alpha _2,\alpha _{2+3}\alpha _2+\alpha _3\},\alpha _2=ϵ_2\delta _1,\alpha _3=\delta _1\delta _2$$
(113)
Using the results on (classical) SCA obtained above for general $`N`$ we may almost immediately complete the (classical) $`N=4`$ SCA. We find that closure is ensured by the following 12 BRST invariant generators which may be characterized as:
Virasoro generator $`L`$ $`h=2`$
supercurrents $`G_{\alpha _2},G_{\alpha _{2+3}},\overline{G}_{\alpha _2},\overline{G}_{\alpha _{2+3}}`$ $`h=3/2`$
$`\text{affine}SL(2)`$ $`\stackrel{~}{E}=K_{\alpha _{2+3};\alpha _2},`$
$`\stackrel{~}{H}=K_{\alpha _{2+3};\alpha _{2+3}}K_{\alpha _2;\alpha _2},`$
$`\stackrel{~}{F}=K_{\alpha _2;\alpha _{2+3}}`$ $`h=1`$
$`\text{affine}U(1)`$ $`U=K_{\alpha _2;\alpha _2}+K_{\alpha _{2+3};\alpha _{2+3}}`$ $`h=1`$
fermions $`\varphi _{\alpha _2}=\mathrm{\Phi }_{\alpha _{2+3};\alpha _{2+3},\alpha _2},`$
$`\varphi _{\alpha _{2+3}}=\mathrm{\Phi }_{\alpha _2;\alpha _2,\alpha _{2+3}}`$ $`h=1/2`$
scalar $`S`$ $`h=0`$ (114)
$`(nS_n=\overline{G}_{\alpha _2,\alpha _{2+3};n}`$ $`h=1)`$
The non-trivial (anti-)commutators are
$`[L_n,A_m]`$ $`=`$ $`((h(A)1)nm)A_{n+m}`$
$`\{G_{\alpha ^{};n+1/2},G_{\beta ^{};m+1/2}\}`$ $`=`$ $`0`$
$`\{G_{\alpha ^{};n+1/2},\overline{G}_{\beta ^{};m1/2}\}`$ $`=`$ $`\delta _{\alpha ^{},\beta ^{}}L_{n+m}+(nm+1)K_{\alpha ^{};\beta ^{};n+m}`$
$`\{\overline{G}_{\alpha _2;n1/2},\overline{G}_{\alpha _{2+3};m1/2}\}`$ $`=`$ $`(nm)(n+m1)S_{n+m1}`$
$`[\stackrel{~}{H}_n,\stackrel{~}{E}_m]`$ $`=`$ $`2\stackrel{~}{E}_{n+m},[\stackrel{~}{H}_n,\stackrel{~}{F}_m]=2\stackrel{~}{F}_{n+m},[\stackrel{~}{E}_n,\stackrel{~}{F}_m]=\stackrel{~}{H}_{n+m}`$
$`[\stackrel{~}{E}_n,G_{\alpha _2;m+1/2}]`$ $`=`$ $`G_{\alpha _{2+3};n+m+1/2},[\stackrel{~}{F}_n,G_{\alpha _{2+3};m+1/2}]=G_{\alpha _2;n+m+1/2}`$
$`[\stackrel{~}{H}_n,G_{\alpha _2;m+1/2}]`$ $`=`$ $`G_{\alpha _2;n+m+1/2},[\stackrel{~}{H}_n,G_{\alpha _{2+3};m+1/2}]=G_{\alpha _{2+3};n+m+1/2}`$
$`[\stackrel{~}{E}_n,\overline{G}_{\alpha _{2+3};m1/2}]`$ $`=`$ $`\overline{G}_{\alpha _2;n+m1/2}n\varphi _{\alpha _2;n+m1/2}`$
$`[\stackrel{~}{H}_n,\overline{G}_{\alpha _2;m1/2}]`$ $`=`$ $`\overline{G}_{\alpha _2;n+m1/2}+n\varphi _{\alpha _2;n+m1/2}`$
$`[\stackrel{~}{H}_n,\overline{G}_{\alpha _{2+3};m1/2}]`$ $`=`$ $`\overline{G}_{\alpha _{2+3};n+m1/2}n\varphi _{\alpha _{2+3};n+m1/2}`$
$`[\stackrel{~}{F}_n,\overline{G}_{\alpha _2;m1/2}]`$ $`=`$ $`\overline{G}_{\alpha _{2+3};n+m1/2}n\varphi _{\alpha _{2+3};n+m1/2}`$
$`[U_n,\overline{G}_{\alpha _2;m1/2}]`$ $`=`$ $`n\varphi _{\alpha _2;n+m1/2}`$
$`[U_n,\overline{G}_{\alpha _{2+3};m1/2}]`$ $`=`$ $`n\varphi _{\alpha _{2+3};n+m1/2}`$
$`[S_n,G_{\alpha _2;m+1/2}]`$ $`=`$ $`\varphi _{\alpha _{2+3};n+m+1/2},[S_n,G_{\alpha _{2+3};m+1/2}]=\varphi _{\alpha _2;n+m+1/2}`$
$`\{G_{\alpha _2;n+1/2},\varphi _{\alpha _2;m1/2}\}`$ $`=`$ $`U_{n+m},\{G_{\alpha _{2+3};n+1/2},\varphi _{\alpha _{2+3};m1/2}\}=U_{n+m}`$
$`\{\overline{G}_{\alpha _2;n1/2},\varphi _{\alpha _{2+3};m1/2}\}`$ $`=`$ $`(n+m1)S_{n+m1}`$
$`\{\overline{G}_{\alpha _{2+3};n1/2},\varphi _{\alpha _2;m1/2}\}`$ $`=`$ $`(n+m1)S_{n+m1}`$
$`[\stackrel{~}{E}_n,\varphi _{\alpha _{2+3};m1/2}]`$ $`=`$ $`\varphi _{\alpha _2;n+m1/2},[\stackrel{~}{F}_n,\varphi _{\alpha _2;m1/2}]=\varphi _{\alpha _{2+3};n+m1/2}`$
$`[\stackrel{~}{H}_n,\varphi _{\alpha _2;m1/2}]`$ $`=`$ $`\varphi _{\alpha _2;n+m1/2},[\stackrel{~}{H}_n,\varphi _{\alpha _{2+3};m1/2}]=\varphi _{\alpha _{2+3};n+m1/2}`$
$`A_m`$ denotes any of the 12 BRST invariant generators listed in (114). We observe that only the derivative of the scalar $`S`$ appears on the right hand sides of (LABEL:N4). Thus, the zero mode of $`S`$ decouples from the algebra. One may verify explicitly that the Jacobi identities are satisfied.
Finally, we note that this center-less $`N=4`$ SCA is of a new and asymmetric form. In particular, it deviates essentially from the small $`N=4`$ SCA announced in Ref. to be obtained by a similar construction. Even though the free field realization of the associated affine $`SL(2|2)`$ current superalgebra used in Ref. is slightly different from ours, one may show that the result in Ref. for the $`N=4`$ SCA is incorrect. One way of reaching this conclusion is to consider the analogue to our (78) and specialize to the case $`n=0`$. In the notation of Ref. , this corresponds to considering the anti-commutator $`\{\overline{G}_{1/2}^1,\overline{G}_{m1/2}^2\}`$ in which case $`\overline{G}_{1/2}^1`$ reduces to $`\frac{dz}{2\pi i}j_{\alpha _1+\alpha _2}`$ (still in the notation of Ref. ). We find that the anti-commutator for generic $`m`$ is non-vanishing in agreement with our result but contrary to the definition of the small $`N=4`$ SCA. Nevertheless, the small $`N=4`$ SCA can be obtained by a construction equivalent to the one employed above. One simply has to replace the original world sheet $`SL(2|2)`$ current superalgebra by an $`SL(2|2)/U(1)`$ current superalgebra, whereby the resulting space-time SCA reduces to the standard small $`N=4`$ SCA. This will be discussed further in Ref. .
## 7 Conclusion
In the present paper a new class of two-dimensional $`N`$ extended SCAs has been discussed. The algebras are induced by free field realizations of affine $`SL(2|N/2)`$ current superalgebras, where $`N`$ is even. In the framework of string theory on $`AdS_3`$ the affine $`SL(2|N/2)`$ current superalgebra resides on the world sheet providing a space-time SCA on the boundary of $`AdS_3`$. The construction generalizes recent work by Ito . The Virasoro generators, the $`N`$ supercurrents, and the generators of an internal $`SL(N/2)U(1)`$ Kac-Moody algebra have all been constructed explicitly. Reducing the considerations to a classical center-less limit has provided additional insight into the structure of the full SCA. BRST invariance has also been addressed. The classical $`N=4`$ SCA is complete and of a new type. In particular, it differs from the small $`N=4`$ SCA. It also illustrates the new and important property of the general construction that it treats the supercurrents asymmetrically.
The results presented here offer (“stringy”) representations of superconformal algebras which are linear in the currents. This suggests that they may be useful when discussing representation theoretical questions, and in the computation of correlation functions. Many other applications may be envisaged.
Several classes of Lie supergroups enjoy decompositions of the bosonic part $`G=SL(2)G^{}`$ as in the case of $`SL(2|N/2)`$. Based on their associated current superalgebras, we anticipate that other classes of SCAs may be constructed along the lines employed in the present paper. This is currently being investigated.
In the classification of CFT with extended symmetries, the construction of SCAs in the present paper presents an alternative to conventional Hamiltonian reduction and otherwise constructed non-linearly extended SCAs . Whole new classes of extended Virasoro algebras seem to be the result of it. There are strong indications that we are even able to produce new and purely bosonic (and linear) extensions of the Virasoro algebra. These will be the subject of a forthcoming publication.
Acknowledgment
The author is grateful to K. Olsen and J.L. Petersen for comments at early stages of this work, and thanks The Niels Bohr Institute, where parts of this work were done, for its kind hospitality.
## Appendix A Lie Superalgebra $`sl(2|M)`$
The root space of the Lie superalgebra $`sl(2|M)`$ in the distinguished representation may be realized in terms of an orthonormal two-dimensional basis $`\{ϵ_1,ϵ_2\}`$ and an orthonormal $`M`$-dimensional basis $`\left\{\delta _u\right\}_{u=1,\mathrm{},M}`$ with metrics
$$ϵ_\iota ϵ_\iota ^{}=\delta _{\iota ,\iota ^{}},\delta _u\delta _u^{}=\delta _{u,u^{}},ϵ_\iota \delta _u=0$$
(116)
The $`\frac{1}{2}(M+1)(M+2)`$ positive roots are then represented as
$`\mathrm{\Delta }_+^0`$ $`=`$ $`\left\{ϵ_1ϵ_2\right\}\left\{\delta _u\delta _v\right|u<v\}`$
$`\mathrm{\Delta }_+^{1+}`$ $`=`$ $`\left\{ϵ_1\delta _u\right|u=1,\mathrm{},M\}`$
$`\mathrm{\Delta }_+^1`$ $`=`$ $`\left\{ϵ_2\delta _u\right|u=1,\mathrm{},M\}`$ (117)
where the $`M+1`$ simple roots $`\alpha _i`$ are
$`\alpha _1`$ $`=`$ $`ϵ_1ϵ_2`$
$`\alpha _2`$ $`=`$ $`ϵ_2\delta _1`$
$`\alpha _{u+2}`$ $`=`$ $`\delta _u\delta _{u+1}`$ (118)
The associated ladder operators $`E_\alpha ,F_\alpha `$, and the Cartan generators $`H_i`$ admit a standard oscillator realization (see e.g. )
$`E_{ϵ_1ϵ_2}=a_1^{}a_2,`$ $`E_{ϵ_\iota \delta _u}=a_\iota ^{}b_u,`$ $`E_{\delta _u\delta _v}=b_u^{}b_v`$
$`F_{ϵ_1ϵ_2}=a_2^{}a_1,`$ $`F_{ϵ_\iota \delta _u}=b_u^{}a_\iota ,`$ $`F_{\delta _u\delta _v}=b_v^{}b_u`$
$`H_1=a_1^{}a_1a_2^{}a_2,`$ $`H_2=a_2^{}a_2+b_1^{}b_1,`$ $`H_{u+2}=b_u^{}b_ub_{u+1}^{}b_{u+1}`$ (119)
where $`a_\iota ^{()}`$ and $`b_u^{()}`$ are fermionic and bosonic oscillators, respectively, satisfying
$$\{a_\iota ,a_\iota ^{}^{}\}=\delta _{\iota ,\iota ^{}},[b_u,b_v^{}]=\delta _{u,v},[b_u^{()},a_\iota ^{()}]=0$$
(120)
It is not possible to design a root string from $`\delta _u\delta _v<0`$ to $`\alpha _1`$ implying that
$$V_{\delta _u\delta _v}^{\alpha _1}=0$$
(121)
This fact is used in deriving (50). However, root strings from $`\alpha ^\pm `$ to $`\alpha _1`$ do exist. They are of the form
$`\alpha ^{}+\mathrm{}+(ϵ_2\delta _u)+\alpha _1`$ $`=`$ $`\alpha _1`$
$`\alpha ^{}+\mathrm{}+(ϵ_1\delta _u)`$ $`=`$ $`\alpha _1`$ (122)
where root strings from $`\alpha ^+`$ are obtained by “inserting” an additional $`\alpha _1`$:
$`\alpha ^++\mathrm{}+\alpha _1+\mathrm{}+(ϵ_2\delta _u)+\alpha _1`$ $`=`$ $`\alpha _1`$
$`\alpha ^++\mathrm{}+(ϵ_2\delta _u)+\alpha _1+\alpha _1`$ $`=`$ $`\alpha _1`$
$`\alpha ^++\mathrm{}+\alpha _1+\mathrm{}+(ϵ_1\delta _u)`$ $`=`$ $`\alpha _1`$
$`\alpha ^++\mathrm{}+(ϵ_1\delta _u)+\alpha _1`$ $`=`$ $`\alpha _1`$ (123)
Possible additions of positive even roots $`\delta _v^{}\delta _v`$ are indicated by “…”. Let us consider the polynomials (20)
$$V_{\alpha ^\pm }^{\alpha _1}=\left[\underset{n0}{}\frac{1}{n!}(C)^n\right]_{\alpha ^\pm }^{\alpha _1}$$
(124)
and compare the two polynomials in order to derive the relation (59). This we do by considering $`\frac{1}{n!}(C)^n`$ in $`V_\alpha ^{}^{\alpha _1}`$ and $`\frac{1}{(n+1)!}(C)^{n+1}`$ in $`V_{\alpha ^+}^{\alpha _1}`$. There are two cases, as such terms involve either $`c^{ϵ_1\delta _u}`$ or $`c^{ϵ_2\delta _u}`$ for some $`u`$, and we will discuss them separately. In the first case the relevant root strings are the lower one in (122) and the two lower ones in (123). Their differences are characterized by the structure constants
$$f_{(ϵ_2\delta _u),(ϵ_1\delta _u)}^{}{}_{}{}^{\alpha _1}=1$$
(125)
and
$`f_{(ϵ_1\delta _v),\alpha _1}^{}{}_{}{}^{(ϵ_2\delta _v)}f_{(ϵ_2\delta _u),(ϵ_1\delta _u)}^{}{}_{}{}^{\alpha _1}`$ $`=`$ $`1`$
$`f_{(ϵ_1\delta _u),ϵ_1\delta _u}^{}{}_{}{}^{j}f_{j,\alpha _1}^{}{}_{}{}^{\alpha _1}`$ $`=`$ $`1`$ (126)
Now, each time the situation (125) occurs in $`\frac{1}{n!}(C)^n`$ in $`V_\alpha ^{}^{\alpha _1}`$, the situations (126) occur $`n`$ times and once, respectively, in $`\frac{1}{(n+1)!}(C)^{n+1}`$ in $`V_{\alpha ^+}^{\alpha _1}`$. This is in accordance with (59). A similar analysis of the terms involving $`c^{ϵ_2\delta _u}`$ leads to the characterizations
$$f_{(ϵ_2\delta _u),ϵ_2\delta _u}^{}{}_{}{}^{j}f_{j,\alpha _1}^{}{}_{}{}^{\alpha _1}=1$$
(127)
and
$`f_{(ϵ_1\delta _v),\alpha _1}^{}{}_{}{}^{(ϵ_2\delta _v)}f_{(ϵ_2\delta _u),ϵ_2\delta _u}^{}{}_{}{}^{j}f_{j,\alpha _1}^{}{}_{}{}^{\alpha _1}`$ $`=`$ $`1`$
$`f_{(ϵ_1\delta _u),ϵ_2\delta _u}^{}{}_{}{}^{\alpha _1}f_{\alpha _1,\alpha _1}^{}{}_{}{}^{1}f_{1,\alpha _1}^{}{}_{}{}^{\alpha _1}`$ $`=`$ $`2`$ (128)
Each time the situation (127) occurs in $`\frac{1}{n!}(C)^n`$ in $`V_\alpha ^{}^{\alpha _1}`$, the situations (128) occur $`n1`$ times and once, respectively, in $`\frac{1}{(n+1)!}(C)^{n+1}`$ in $`V_{\alpha ^+}^{\alpha _1}`$. However, as the last situation contributes with a factor $`2`$, the terms involving $`c^{ϵ_2\delta _u}`$ also agree with the relation (59), which is thereby proven.
### A.1 Polynomial Relations
Using that the polynomials $`V`$ enter in a differential operator realization of a Lie superalgebra leads to the polynomial relations
$$f_{a,b}^{}{}_{}{}^{c}V_c^\alpha =V_a^\nu _\nu V_b^\alpha (1)^{p(a)p(b)}V_b^\nu _\nu V_a^\alpha $$
(129)
as discussed in Ref. . The parity $`p(a)`$ of the index $`a`$ is defined as 1 for $`a`$ an odd root and 0 otherwise. Particularly useful are the following relations:
$`V_{\alpha _1}^\nu _\nu V_\alpha ^{}^{\alpha _1}`$ $`=`$ $`0`$
$`V_1^\nu _\nu V_\alpha ^{}^{\alpha _1}`$ $`=`$ $`V_\alpha ^{}^{\alpha _1}`$
$`V_{\alpha _1}^\nu _\nu V_\alpha ^{}^{\alpha _1}`$ $`=`$ $`\gamma V_\alpha ^{}^{\alpha _1}`$ (130)
$`V_{\alpha ^+}^\mu _\mu V_\beta ^{}^{\alpha _1}`$ $`=`$ $`\delta _{\alpha ^{},\beta ^{}}`$
$`V_\alpha ^{}^\mu _\mu V_\beta ^{}^{\alpha _1}`$ $`=`$ $`\gamma \delta _{\alpha ^{},\beta ^{}}`$
$`V_{\alpha ^+}^\mu _\mu V_\beta ^{}^{\alpha _1}`$ $`=`$ $`V_\alpha ^{}^{\alpha _1}V_\beta ^{}^{\alpha _1}`$
$`V_\alpha ^{}^\mu _\mu V_\beta ^{}^{\alpha _1}`$ $`=`$ $`0`$ (131)
$`f_{\alpha ^{},\beta ^{}}^{}{}_{}{}^{c}V_c^{\alpha _1}`$ $`=`$ $`\gamma \delta _{\alpha ^{},\beta ^{}}`$
$`f_{\alpha ^{},\beta ^{}}^{}{}_{}{}^{c}V_c^\mu _\mu V_\nu ^{}^{\alpha _1}`$ $`=`$ $`\delta _{\alpha ^{},\nu ^{}}V_\beta ^{}^{\alpha _1}`$ (132)
$`f_{\alpha ^+,\beta ^{}}^{}{}_{}{}^{c}J_c`$ $`=`$ $`\delta _{\alpha ^{},\beta ^{}}E_1`$
$`f_{\alpha ^+,\beta ^+}^{}{}_{}{}^{c}J_c`$ $`=`$ $`\delta _{\alpha ^{},\beta ^{}}H_1+f_{\alpha ^{},\beta ^{}}^{}{}_{}{}^{c}J_c`$
$`f_{\alpha ^{},\beta ^+}^{}{}_{}{}^{c}J_c`$ $`=`$ $`\delta _{\alpha ^{},\beta ^{}}F_1`$ (133)
$`f_{\alpha ^{},\beta ^{}}^{}{}_{}{}^{c}f_{c,\lambda ^{}}^{}{}_{}{}^{d}J_d`$ $`=`$ $`\delta _{\alpha ^{},\beta ^{}}J_\lambda ^{}\delta _{\beta ^{},\lambda ^{}}J_\alpha ^{}`$
$`f_{\alpha ^{},\beta ^{}}^{}{}_{}{}^{c}f_{c,\lambda ^+}^{}{}_{}{}^{d}J_d`$ $`=`$ $`\delta _{\beta ^{},\lambda ^{}}J_{\alpha ^+}`$
$`f_{\alpha ^{},\beta ^{}}^{}{}_{}{}^{c}f_{c,\lambda ^{}}^{}{}_{}{}^{d}J_d`$ $`=`$ $`\delta _{\alpha ^{},\lambda ^{}}J_\beta ^{}\delta _{\alpha ^{},\beta ^{}}J_\lambda ^{}`$
$`f_{\alpha ^{},\beta ^{}}^{}{}_{}{}^{c}f_{c,\lambda ^+}^{}{}_{}{}^{d}J_d`$ $`=`$ $`\delta _{\alpha ^{},\lambda ^{}}J_{\beta ^+}`$
$`f_{\alpha ^{},\beta ^{}}^{}{}_{}{}^{c}f_{c,\alpha _1}^{}{}_{}{}^{d}J_d`$ $`=`$ $`\delta _{\alpha ^{},\beta ^{}}E_1`$
$`f_{\alpha ^{},\beta ^{}}^{}{}_{}{}^{c}f_{c,1}^{}{}_{}{}^{d}J_d`$ $`=`$ $`0`$
$`f_{\alpha ^{},\beta ^{}}^{}{}_{}{}^{c}f_{c,\alpha _1}^{}{}_{}{}^{d}J_d`$ $`=`$ $`\delta _{\alpha ^{},\beta ^{}}F_1`$ (134)
$$f_{\alpha ^{},\beta ^{}}^{}{}_{}{}^{c}f_{\mu ^{},\nu ^{}}^{}{}_{}{}^{d}f_{c,d}^{}{}_{}{}^{e}J_e=\delta _{\alpha ^{},\nu ^{}}f_{\mu ^{},\beta ^{}}^{}{}_{}{}^{c}J_c\delta _{\mu ^{},\beta ^{}}f_{\alpha ^{},\nu ^{}}^{}{}_{}{}^{c}J_c$$
(135)
In deriving some of these relations we have made use of the explicit oscillator realization (119). |
warning/0002/cond-mat0002058.html | ar5iv | text | # A Mean Atom Trajectory Model for Monatomic Liquids
## 1 Introduction
The motion of atoms in a liquid can be divided into two constituent parts: (a) oscillation in a valley of the liquid’s many-body potential and (b) transits between many-body valleys. The latter process is responsible for self-diffusion. As we showed in a previous paper , which drew upon earlier results of Clements and Wallace , the former motion can be modeled very precisely under the assumption that the valleys are nearly harmonic, with the majority of valleys (the random valleys) sharing a common spectrum of frequencies. Specifically, a purely harmonic model provides an extremely accurate formula for $`\widehat{Z}(t)`$, the normalized velocity autocorrelation function, in the nondiffusing regime. Here we will use this picture and our previous work to justify a “mean atom trajectory” model, a single-atom model that approximates the behavior of an average atom in the liquid and correctly reproduces its nondiffusing behavior. Then we will introduce a simple intuitive account of the transit process that allows us to extend the model to the self-diffusing regime. In Sec. 2 we develop this model and explain how it is used to calculate $`\widehat{Z}(t)`$ for a diffusing liquid. Then we fit the model to MD simulations at various temperatures in Sec. 3, and we comment on the quality of the results. In Sec. 4 we compare the present mean atom trajectory model with previous work based on Instantaneous Normal Modes (INM), and with an earlier independent atom model , and we summarize our conclusions.
## 2 The Mean Atom Trajectory Model
### 2.1 General comments
To form an appropriate basis for our model, we must begin with some initial reasonable approximations about the nature of the valleys and transits in the real liquid. As mentioned in the Introduction, available evidence suggests that the valleys are nearly harmonic, and we will continue to assume that here. Further, we will assume that transits between valleys occur instantaneously and are local in character; that is, each transit involves only a few neighboring atoms.
Now for the nondiffusing supercooled liquid, there are no transits, and as we’ve seen the system’s motion is accurately expressed in terms of the harmonic normal modes. But it is also legitimate to consider this motion from the point of view of a single atom, and when we do so we see that each atom moves along a complicated trajectory within a single-particle potential well which fluctuates because of the motion of neighboring atoms, but whose center is fixed in space. An important timescale for this motion is the single-atom mean vibrational period $`\tau `$, which we take as $`2\pi /\omega _{\mathrm{rms}}`$, where $`\omega _{\mathrm{rms}}`$ is the rms frequency of the set of normal modes. Let us now follow the single-atom description as the temperature is increased. Once the glass transition is passed, the atom will begin to make transits from one single-particle well to another , and even before the melting temperature is reached, the transit rate will be on the order of one per mean vibrational period. Since each atom has approximately ten neighbors, roughly ten transits will occur in its immediate vicinity every period, changing the set of normal mode eigenvectors each time; thus we conclude that a decomposition of the motion into normal modes will not be useful when the liquid is diffusing, so we have no choice but to follow a single-atom description for the liquid state . That being the case, we shall start from the beginning with a single-atom description, in order to construct a unified model for diffusing and nondiffusing motion alike.
### 2.2 Nondiffusing regime
Our starting point will be in the nondiffusing regime, where a normal mode analysis is still valid. As shown in , the $`i`$th coordinate of the $`K`$th atom in an $`N`$-body harmonic valley can be written
$$u_{Ki}(t)=\underset{\lambda }{}w_{Ki,\lambda }a_\lambda \mathrm{sin}(\omega _\lambda t+\alpha _\lambda ),$$
(1)
where the $`w_{Ki,\lambda }`$ form a $`3N\times 3N`$ orthogonal matrix, the $`\omega _\lambda `$ are the frequencies of the normal modes, and the $`a_\lambda `$ are the amplitudes of the modes. Three of the modes have zero frequency and correspond to center of mass motion; here we demand that the center of mass is stationary, so the system has only $`3N3`$ independent degrees of freedom, the zero frequency modes are absent, and the sum over $`\lambda `$ runs from 1 to $`3N3`$. To make this an equation for a “mean” atom, we first drop the index $`K`$:
$$u_i(t)=\underset{\lambda }{}w_{i\lambda }a_\lambda \mathrm{sin}(\omega _\lambda t+\alpha _\lambda ).$$
(2)
Now we must reinterpret the $`w_{i\lambda }`$ since they no longer form an orthogonal (or even square) matrix. Let
$$𝒘_\lambda =w_{1\lambda }\widehat{𝒙}+w_{2\lambda }\widehat{𝒚}+w_{3\lambda }\widehat{𝒛}$$
(3)
so
$$𝒖(t)=\underset{\lambda }{}𝒘_\lambda a_\lambda \mathrm{sin}(\omega _\lambda t+\alpha _\lambda ).$$
(4)
We will ultimately consider situations in which the well center is allowed to move, so let $`𝒓`$$`(t)`$ be the atom’s position, $`𝑹`$ be the location of the center of the well, and $`𝒖`$$`(t)`$ be the atom’s displacement from the well center; then
$$𝒓(t)=𝑹+\underset{\lambda }{}𝒘_\lambda a_\lambda \mathrm{sin}(\omega _\lambda t+\alpha _\lambda )$$
(5)
with velocity
$$𝒗(t)=\underset{\lambda }{}𝒘_\lambda a_\lambda \omega _\lambda \mathrm{cos}(\omega _\lambda t+\alpha _\lambda ).$$
(6)
We now have the basic formula, but we must decide how to assign values to the $`𝒘`$<sub>λ</sub>, $`a_\lambda `$, and $`\alpha _\lambda `$. Let’s do this by calculating $`Z(t)`$, the velocity autocorrelation function, in this model and comparing it to the harmonic result derived in .
$`Z(t)`$ $`=`$ $`{\displaystyle \frac{1}{3}}𝒗(t)𝒗(0)`$ (7)
$`=`$ $`{\displaystyle \frac{1}{3}}{\displaystyle \underset{\lambda \lambda ^{}}{}}𝒘_\lambda 𝒘_\lambda ^{}a_\lambda a_\lambda ^{}\omega _\lambda \omega _\lambda ^{}\mathrm{cos}(\omega _\lambda t+\alpha _\lambda )\mathrm{cos}(\alpha _\lambda ^{})`$
Let’s assign the $`\alpha _\lambda `$ randomly and average over each $`\alpha _\lambda `$ separately; then if $`\lambda \lambda ^{}`$,
$$\mathrm{cos}(\omega _\lambda t+\alpha _\lambda )\mathrm{cos}(\alpha _\lambda ^{})=\mathrm{cos}(\omega _\lambda t+\alpha _\lambda )\mathrm{cos}(\alpha _\lambda ^{})=0$$
(8)
but if $`\lambda =\lambda ^{}`$, then
$$\mathrm{cos}(\omega _\lambda t+\alpha _\lambda )\mathrm{cos}(\alpha _\lambda ^{})=\mathrm{cos}(\omega _\lambda t+\alpha _\lambda )\mathrm{cos}(\alpha _\lambda )=\frac{1}{2}\mathrm{cos}(\omega _\lambda t).$$
(9)
Thus the $`\lambda \lambda ^{}`$ terms in Eq. (7) are eliminated and we have
$$Z(t)=\frac{1}{6}\underset{\lambda }{}|𝒘_\lambda |^2a_\lambda ^2\omega _\lambda ^2\mathrm{cos}(\omega _\lambda t).$$
(10)
Note that the formula for $`Z(t)`$ in the harmonic theory also lacks off-diagonal terms , but for a different reason: The orthogonality of the matrix $`w_{Ki,\lambda }`$ removes them in that case. Now let us make the assignment
$$𝒘_\lambda =\frac{1}{\sqrt{N1}}\widehat{𝒘}_\lambda $$
(11)
where $`\widehat{𝒘}_\lambda `$ is a randomly chosen unit vector; then
$$Z(t)=\frac{1}{6N6}\underset{\lambda }{}a_\lambda ^2\omega _\lambda ^2\mathrm{cos}(\omega _\lambda t).$$
(12)
This is the same expression for $`Z(t)`$ that one derives in the harmonic model (see Eq. (10) of ) for any distribution of normal mode amplitudes $`a_\lambda `$, not just the thermal equilibrium distribution. Thus our model with these choices of $`\alpha _\lambda `$ and $`𝒘`$<sub>λ</sub> correctly reproduces $`Z(t)`$ for any equilibrium or nonequilibrium ensemble in the harmonic theory. To recover the equilibrium result, we make the final substitution
$$a_\lambda =\sqrt{\frac{2kT}{M\omega _\lambda ^2}}$$
(13)
with the result
$$Z(t)=\frac{1}{3N3}\frac{kT}{M}\underset{\lambda }{}\mathrm{cos}(\omega _\lambda t),$$
(14)
which is Eq. (13) from . Thus our model assumes that the motion of a mean nondiffusing atom in thermal equilibrium at temperature $`T`$ is given by
$$𝒓(t)=𝑹+\frac{1}{\sqrt{N1}}\sqrt{\frac{2kT}{M}}\underset{\lambda }{}\widehat{𝒘}_\lambda \omega _\lambda ^1\mathrm{sin}(\omega _\lambda t+\alpha _\lambda )$$
(15)
where the phases $`\alpha _\lambda `$ and unit vectors $`\widehat{𝒘}_\lambda `$ are randomly chosen.
By construction this model gets the same result for $`Z(t)`$, and thus $`v^2`$, as the harmonic model; does it correctly reproduce any other functions? We can check by calculating the analog of $`Z(t)`$ for positions, $`𝒖(t)𝒖(0)`$. Using Eq. (1) for $`u_{Ki}(t)`$ and calculating for the harmonic model as in , and using Eq. (2) for $`𝒖`$$`(t)`$ and calculating as above, one finds in both cases that
$$𝒖(t)𝒖(0)=\frac{1}{2N2}\underset{\lambda }{}a_\lambda ^2\mathrm{cos}(\omega _\lambda t).$$
(16)
Again the substitution $`a_\lambda =\sqrt{2kT/M\omega _\lambda ^2}`$ recovers the correct thermal equilibrium result. So this model correctly reproduces $`𝒖(t)𝒖(0)`$ and $`u^2`$ as well. Notice that because we have a closed form for $`𝒖`$$`(t)`$, all of the correlation functions calculated so far also have a closed form. Once we introduce transits, this will no longer be the case.
While this model compares well with the harmonic model, one might wonder how well it compares with molecular dynamics (MD) results. In Fig. 1 a graph of $`v^2`$ as a function of $`t`$ for the model in equilibrium at 6.69 K is compared to a randomly chosen particle from an MD run
of liquid Na also at 6.69 K, a temperature at which it is known the sample is nondiffusing (see for details). Note that in both graphs $`v^2`$ has approximately the same amplitude, and the gaps between peaks are roughly the same size, indicating oscillations at frequencies in the same ranges. Thus not only functions of the motion but the motion itself shows strong qualitative agreement with MD in the nondiffusing regime.
The MD system with which we will be comparing our model has $`N=500`$ particles, so it has 1497 normal mode frequencies. (Remember that the zero frequency modes are removed at the start.) When we introduce transits, we will have to evaluate $`Z(t)`$ numerically, and for speed of computation we would like to use only a representative subset of the normal mode frequencies, say 75 instead of all 1497. The original set of frequencies is determined as described in . We decided which subset to use by calculating three moments of the full frequency distribution defined below:
$`\omega _2`$ $`=`$ $`\left[{\displaystyle \frac{5}{3}}\omega _\lambda ^2\right]^{1/2}`$
$`\omega _0`$ $`=`$ $`\mathrm{exp}\mathrm{ln}(\omega _\lambda )`$
$`\omega _2`$ $`=`$ $`\left[{\displaystyle \frac{1}{3}}\omega _\lambda ^2\right]^{1/2}.`$ (17)
Note that by this definition $`\omega _{\mathrm{rms}}=\sqrt{3/5}\omega _2`$. We then calculated the same three moments for several sets of 75 frequencies evenly spaced throughout the full set, and we chose the set that best fit the moments to use in further calculations. This process is somewhat subjective, because skewing the sample in favor of lower frequencies improves the accuracy of $`\omega _2`$ but reduces the accuracy of $`\omega _2`$, and the opposite is true if one skews in favor of high frequencies. The table below shows the values of the moments for the full set and the reduced set of 75 that we ultimately chose. The frequencies are in units of $`\delta t^1`$, where $`\delta t=1.4\times 10^{15}`$ s is the timestep of our MD simulations .
| | Full | Reduced |
| --- | --- | --- |
| $`\omega _2`$ | 0.02826 | 0.02824 |
| $`\omega _0`$ | 0.01807 | 0.01803 |
| $`\omega _2`$ | 0.02101 | 0.02111 |
To use the reduced set we must rewrite our formulas for $`𝒓`$$`(t)`$ and $`𝒗`$$`(t)`$ slightly. Let $`\mathrm{\Lambda }`$ be the total number of frequencies; then $`\mathrm{\Lambda }=3N3`$ for the full set so $`𝒓`$$`(t)`$ and $`𝒗`$$`(t)`$ become
$`𝒓(t)`$ $`=`$ $`𝑹+\sqrt{{\displaystyle \frac{3}{\mathrm{\Lambda }}}}\sqrt{{\displaystyle \frac{2kT}{M}}}{\displaystyle \underset{\lambda =1}{\overset{\mathrm{\Lambda }}{}}}\widehat{𝒘}_\lambda \omega _\lambda ^1\mathrm{sin}(\omega _\lambda t+\alpha _\lambda )`$
$`𝒗(t)`$ $`=`$ $`\sqrt{{\displaystyle \frac{3}{\mathrm{\Lambda }}}}\sqrt{{\displaystyle \frac{2kT}{M}}}{\displaystyle \underset{\lambda =1}{\overset{\mathrm{\Lambda }}{}}}\widehat{𝒘}_\lambda \mathrm{cos}(\omega _\lambda t+\alpha _\lambda ).`$ (18)
This form is also correct for the reduced set of frequencies $`\omega _\lambda `$, so this is the form we will use. As another check on the accuracy of our results with only 75 frequencies, we recalculated the original $`\mathrm{cos}(\omega _\lambda t)`$ expression for $`\widehat{Z}(t)`$ using both the full and reduced sets; their disagreement is at most 0.01 out to time $`1000\delta t`$. However, the discrepancy grows beyond that time, as the expression using the reduced set begins to experience revivals. Hence we will not consider $`\widehat{Z}(t)`$ beyond that point.
### 2.3 Diffusing regime
To incorporate diffusion into the mean atom trajectory model, we rely on Wallace’s notion of a single-particle transit , a nearly instantaneous transition from one well to another. As discussed in , we expect transits to be governed not by thermal activation (having enough energy to escape a fixed well) but by correlations (neighbors must be positioned properly for a low-potential path to open between wells). We implement this property by having transits occur at a temperature dependent rate $`\nu (T)`$, so in a small time interval $`\mathrm{\Delta }t`$ the probability of a single transit is $`\nu \mathrm{\Delta }t`$.
We model the transit process itself by assuming it occurs instantaneously in the forward direction; from this we can determine the parameters $`𝑹`$, $`\widehat{𝒘}_\lambda `$, and $`\alpha _\lambda `$ appearing in Eq. (18) after the transit in terms of the same quantities before the transit. Since the process is instantaneous, both $`𝒓`$$`(t)`$ and $`𝒗`$$`(t)`$ are the same before and afterwards. Let $`𝑹`$<sup>before</sup> and $`𝑹`$<sup>after</sup> be the well centers before and after the transit and let $`𝒖`$<sup>before</sup> and $`𝒖`$<sup>after</sup> be the corresponding displacements from the well centers. Then $`𝒓^{\mathrm{before}}=𝒓^{\mathrm{after}}`$ implies
$$𝑹^{\mathrm{before}}+𝒖^{\mathrm{before}}=𝑹^{\mathrm{after}}+𝒖^{\mathrm{after}}.$$
(19)
To transit forward, we assume the center of the new well lies along the line between the old well center and the atom, but it lies on the opposite side of the atom from the old well center an equal distance away. This implies $`𝒖`$$`{}_{}{}^{\mathrm{after}}=𝒖^{\mathrm{before}}`$, so
$$𝑹^{\mathrm{before}}+𝒖^{\mathrm{before}}=𝑹^{\mathrm{after}}𝒖^{\mathrm{before}}$$
(20)
with the result
$$𝑹^{\mathrm{after}}=𝑹^{\mathrm{before}}+2𝒖^{\mathrm{before}}.$$
(21)
This determines the new well center in terms of the coordinates before the transit. As for the unit vectors $`\widehat{𝒘}_\lambda `$, since they are randomly generated and play no role in calculating $`\widehat{Z}(t)`$, we have decided to leave them unchanged by transits. Finally, we have the phases $`\alpha _\lambda `$. We must use these to implement the relations
$$𝒖^{\mathrm{after}}=𝒖^{\mathrm{before}},𝒗^{\mathrm{after}}=𝒗^{\mathrm{before}}$$
(22)
which we have assumed above. Since $`𝒖`$$`(t)`$ is a sum of sines and $`𝒗`$$`(t)`$ a sum of cosines, the simplest way to change the sign of $`𝒖`$$`(t)`$ while preserving that of $`𝒗`$$`(t)`$ is to reverse the signs of the arguments $`(\omega _\lambda t+\alpha _\lambda )`$ in Eq. (18). Let the transit occur at time $`t_0`$; then
$$\omega _\lambda t_0+\alpha _\lambda ^{\mathrm{after}}=(\omega _\lambda t_0+\alpha _\lambda ^{\mathrm{before}})$$
(23)
so
$$\alpha _\lambda ^{\mathrm{after}}=2\omega _\lambda t_0\alpha _\lambda ^{\mathrm{before}}.$$
(24)
Thus, in this model a transit is implemented at time $`t_0`$ by leaving the $`\widehat{𝒘}_\lambda `$ alone and making the substitutions
$`𝑹`$ $``$ $`𝑹+2𝒖(t_0)`$
$`\alpha _\lambda `$ $``$ $`2\omega _\lambda t_0\alpha _\lambda .`$ (25)
This conserves $`𝒓`$, reverses the sign of $`𝒖`$, and conserves $`𝒗`$.
Now our mean atom trajectory model consists of nondiffusive motion between transits as given by Eq. (18), with a given probability in each small time interval that $`𝑹`$ and the phases $`\alpha _\lambda `$ will be replaced with new values as determined in Eq. (25). The addition of transits means that we no longer have closed form expressions for $`𝒓`$$`(t)`$ and $`𝒗`$$`(t)`$ for all times, so we have no closed form expression for $`\widehat{Z}(t)`$; but this model can be implemented easily on a computer in a manner analogous to an MD calculation, and in that way we can calculate autocorrelation functions. We turn to that calculation next.
### 2.4 Evaluating $`\widehat{Z}(t)`$ in the diffusing regime
To calculate $`\widehat{Z}(t)`$, we select a value for the rate $`\nu `$, generate a random set of $`\widehat{𝒘}_\lambda `$ and $`\alpha _\lambda `$, and use Eq. (18) to calculate $`𝒓`$$`(t)`$ and $`𝒗`$$`(t)`$ from $`t=0`$ to $`t=t_{\mathrm{max}}`$ in increments of $`\delta t`$, where the criterion for choosing $`t_{\mathrm{max}}`$ is discussed below and $`\delta t`$ is the timestep used in our MD simulations (defined in Subsection 2.2). At each timestep, we check to see if a transit occurs, and if so we implement Eq. (25) and continue with the new $`𝑹`$ and $`\alpha _\lambda `$. We then calculate $`\widehat{Z}(t)`$ using the formula
$$Z(t)=\frac{1}{3(t_{\mathrm{max}}t)+3}\underset{t^{}=0}{\overset{t_{\mathrm{max}}t}{}}𝒗(t+t^{})𝒗(t^{})$$
(26)
and normalizing. This equation is a modified form of the expression used to calculate $`Z(t)`$ in MD; notice that the average over $`t^{}`$ has the same effect as averaging separately over each phase $`\alpha _\lambda `$ that appears in the velocity vectors. Just as in MD, we want to average over a large data set, so we require $`t_{\mathrm{max}}>>t`$; we have chosen $`t_{\mathrm{max}}=20`$ million timesteps and we calculate $`\widehat{Z}(t)`$ only to $`t=1000`$ timesteps. We estimate the total error from using only a subset of all 1497 frequencies and the finite size of the data set to be at most 0.01; in particular, when $`\nu =0`$ the calculation converges to the closed form result $`\mathrm{cos}(\omega _\lambda t)`$ to this accuracy.
## 3 Comparison with MD
The MD setup with which we compared the predictions of this model is the one described in ; $`N=500`$ atoms of Na move under the influence of a highly realistic pair potential with the timestep $`\delta t=1.4\times 10^{15}`$ given in Subsection 2.2. We performed equilibrium runs of the system at 216.3 K, 309.7 K, 425.0 K, 664.7 K, and 1022.0 K, all temperatures at which the system is diffusing. Since $`T_m=371.0`$ K for Na at this density, our simulations range from the supercooled regime to nearly three times the melting temperature. We then ran the model for various values of $`\nu `$, adjusting until the model matched the value of the first minimum of $`\widehat{Z}(t)`$ at each temperature. The values of $`\nu `$ that we fit for all temperatures are given below, and the resulting $`\widehat{Z}(t)`$ for each $`\nu `$ is compared to the corresponding MD result in Figs. 2 through 6. Here $`\nu `$ is expressed in units of $`\tau ^1`$ where $`\tau =2\pi /\omega _{\mathrm{rms}}`$ is the single-atom mean vibrational period defined in Subsection 2.1.
| $`T`$ (K) | $`\nu (\tau ^1)`$ |
| --- | --- |
| 216.3 | 0.35018 |
| 309.7 | 0.60276 |
| 425.0 | 0.83985 |
| 664.7 | 1.24858 |
| 1022.0 | 1.68774 |
Notice that in all cases $`\nu `$ is of the same order of magnitude as $`\tau ^1`$, indicating roughly one transit per mean vibrational period, as mentioned in Subsection 2.1, and as predicted in .
The most obvious trend exhibited by $`\widehat{Z}(t)`$ from the five MD runs is that its first minimum is rising with increasing $`T`$; as we mention below, this is the primary reason for the increasing diffusion coefficient $`D`$. Note that the model is able to reproduce this most important feature quite satisfactorily. In fact, all five fits of the model to the MD results capture their essential features, but we do see systematic trends in the discrepancies. First, note that the location of the first minimum barely changes at all in the model as $`\nu `$ is raised, but in MD the first minimum moves steadily to earlier times as the temperature rises. The first minimum occurs at a time roughly equal to half of the mean vibrational period ($`\tau =287\delta t`$ in this system), so the steady drift backward suggests that the MD system is sampling a higher range of frequencies at higher $`T`$. Also, for the three lowest temperatures the model tends to overshoot the MD result in the vicinity of the first two maxima after the origin, and at the highest two temperatures this overshoot is accompanied by a positive tail that is slightly higher than the (still somewhat long) tail predicted by MD. These overshoots should clearly affect the diffusion coefficient $`D`$, which is the integral of $`Z(t)`$. To check this, we calculated the reduced diffusion coefficient $`\widehat{D}`$, the integral of $`\widehat{Z}(t)`$, which is related to $`D`$ by
$$D=\frac{kT}{M}\widehat{D}.$$
(27)
The results are compared to the values of $`\widehat{D}`$ calculated from the MD runs in Fig. 7. The results from the two nondiffusing runs discussed in are also included.
In all of the diffusing cases, the model overestimates $`\widehat{D}`$ by roughly the same amount, which we take to be the effect of the overshoots at the first two maxima. At the higher temperatures the discrepancy is also higher, presumably due to the model’s long tail.
It is interesting to note that this MD system produces results that agree very closely with experiment: For example, the MD predicts that at $`T=425.0`$ K and $`\rho =0.925`$ g/cm<sup>3</sup>, $`D=6.40\times 10^5`$ cm<sup>2</sup>/s, while experiments by Larsson, Roxbergh, and Lodding find that at $`T=425.0`$ K and $`\rho =0.915`$ g/cm<sup>3</sup>, $`D=6.020\times 10^5`$ cm<sup>2</sup>/s. Hence our agreement with MD results genuinely reflects agreement with properties of real liquid Na.
## 4 Conclusions
We have presented a single-atom model of a monatomic liquid that provides a unified account of diffusing and nondiffusing behavior. The nondiffusing motion is modeled as a sum of oscillations at the normal mode frequencies (Eq. (18)), simulating the trajectory of an average atom in a complicated single-body potential well that fluctuates due to the motion of its neighbors. Self-diffusion is accounted for in terms of instantaneous transits between wells, which occur at a temperature-dependent rate $`\nu `$. Since this model gives a simple and straightforward account of the motion itself, it can easily be used to calculate any single-atom correlation function one wishes; here we have focussed on the velocity autocorrelation function. It is interesting to note that in this model the velocity correlations persist through a transit, instead of being washed out entirely by the transit process; we will return to this point below. The relaxation of correlations seen by the decay of $`\widehat{Z}(t)`$ arises here from two distinct processes: Dephasing as a result of the large number of frequencies in the single-well motion, and transits between wells. The dephasing effect produces relaxation but not diffusion: It causes $`\widehat{Z}(t)`$ to decay but its integral remains zero. On the other hand, transits certainly contribute to relaxation (see , where they provide the only relaxation mechanism), but in addition they raise the first minimum of $`\widehat{Z}(t)`$ substantially, increasing its integral and providing a nonzero $`D`$.
The comparison of this model to MD results is generally quite positive, particularly for a one-parameter model; the two calculations of $`\widehat{Z}(t)`$ agree strongly over 1000 timesteps. In addition, the match is encouraging over a very large range of temperatures, from essentially 0 K to $`3T_m`$. The most noticeable discrepancies are the backward drift in the location of the first minimum of $`\widehat{Z}(t)`$, which is present in MD but not the model, and the tendency of the model to exaggerate certain characteristics of the MD results (the maxima at intermediate times and the high-$`T`$ positive tail). This latter effect is responsible for the model’s overestimate of the diffusion coefficient, though we hasten to add that the model $`\widehat{D}`$ is still in satisfactory agreement with MD results, especially for the liquid at $`TT_m`$ (see Fig. 7).
As in , it is useful to compare this model and the accompanying results to the work others have done using the formalism of Instantaneous Normal Modes (INM) and similar methods. Previously, we discussed the advantages of our general approach and the superior quality of its results when applied to nondiffusing states; here we will consider matters relating explicitly to mechanisms of diffusion. As is noted explicitly by Vallauri and Bermejo , the account of INM by Stratt (see, for example, ) does not consider diffusion at all; their $`Z(t)`$ is essentially a sum over cosines, and as such it integrates to zero. This is understandable, because as Stratt et al. repeatedly emphasize, their approximation is valid only for very short times, so they are not attempting to model effects with longer timescales. They compare their INM results with MD calculations of $`\widehat{Z}(t)`$ for states of an LJ system ranging from a moderately supercooled liquid to well above the melting temperature, and our fits are of roughly the same quality or better in all cases. Authors who do attempt to model diffusion usually follow the path suggested initially by Zwanzig , who thought of the liquid’s phase space as divided into “cells” in which each atom spends its time before finding a saddle point in the potential and jumping from one cell to the next. He imagined as a first approximation that the jumps destroy all correlations between the cells; since atoms are jumping all the time, he suggested that the net result was to multiply the nondiffusing form of $`Z(t)`$ by a damping factor $`\mathrm{exp}(t/\tau )`$ for some timescale $`\tau `$ representing the lifetime of a stay in a single cell. Notice that Zwanzig provided no dynamical model of the jumping process itself. This suggestion has been developed and transformed extensively by Madan, Keyes, and Seeley , who take a general Zwanzig-like functional form for $`Z(t)`$ and use a combination of heuristic arguments and constraints on its moments to specify its dependence on a “hopping rate” $`\omega _v`$, the analog of $`\tau ^1`$ for Zwanzig, which they then extract from the unstable lobe of the INM spectrum. Although we cannot be entirely sure, as indicated in , we think it most likely that their simulations of LJ Ar at 80 K, 120 K, and 150 K consist of states comparable to our Figs. 2 and 3, and again we would argue that our fits are somewhat better. Finally, Cao and Voth also approach diffusion by means of a damping factor, and they consider factors of two different types, each of which contains parameters that can be determined from other calculated or experimental quantities. Their matches with MD are actually quite good, but again ours are of at least comparable quality.
Having claimed that our matches to MD simulations are as good as or better than all of the others we have surveyed, let us emphasize a fundamental difference between our approaches to diffusion: We provide an account of the process of transiting from well to well, so we have a model of the actual motion of a mean atom in space that is valid to arbitrary times, and given this model we can calculate the effect of transits on correlations. All other approaches we know of begin with an account of the motion valid only for very short times, calculate $`Z(t)`$ from this motion, and then try to model the effects of diffusion on $`Z(t)`$ directly, using parameters that are thought to be characteristic of jumps between wells. In the diffusing regime, no one else we have seen actually calculates $`𝒗(t)𝒗(0)`$ to find $`Z(t)`$, as we do. In the process, we find that some of the assumptions made by others, in particular Zwanzig’s hypothesis that jumps between wells simply erase correlations, are not true. This approach is already yielding insights into the actual motion of an atom undergoing a transit.
Finally, we would like to compare the present mean atom trajectory model with an earlier independent atom model by Wallace . In developing the independent atom model, two arguments were made: (a) the high rate of transits in the liquid state shows the need to abandon the normal mode description of motion, and instead picture the motion of a single atom among a set of fluctuating wells, and (b) the leading approximation to a fluctuating well is its smooth time average well. Accordingly, the independent atom oscillates with frequency $`\omega `$ in a smooth isotropic well, and it transits with probability $`\mu `$ to an adjacent identical well at each turning point . What we have now learned by considering low temperatures, and especially the nondiffusing states of , is that the MD system exhibits a velocity decorrelation process which results from the presence of many frequencies in the single-particle motion, making assumption (b) less reasonable. In the current model these many frequencies are retained, and they are intended to represent the strong fluctuations in each single-particle well; in this way the current model makes an important improvement over the independent atom model, where the well fluctuations were averaged out. Beyond this difference, the two models contain similar but not identical treatments of transits, whose rate increases with increasing temperature, and which produce self-diffusion. The less detailed but simpler independent atom model has proven useful in a description of the glass transition , and the corresponding transit parameter has been used to relate shear viscosity and self-diffusion in liquid metals .
The next logical step in this work, given the results so far, is to extend the mean atom trajectory model and apply it to calculations of more complicated correlation functions of the liquid’s motion. |
warning/0002/cond-mat0002286.html | ar5iv | text | # Ultrafast optical nonlinearity in quasi-one-dimensional Mott-insulator Sr₂CuO₃
## Abstract
We report strong instantaneous photoinduced absorption in the quasi-one-dimensional Mott insulator $`\mathrm{Sr}_2\mathrm{CuO}_3`$ in the IR spectral region. The observed photoinduced absorption is to an even-parity two-photon state that occurs immediately above the absorption edge. Theoretical calculation based on a two-band extended Hubbard model explains the experimental features and indicates that the strong two-photon absorption is due to a very large dipole-coupling between nearly degenerate one- and two-photon states. Room temperature picosecond recovery of the optical transparency suggests the strong potential of $`\mathrm{Sr}_2\mathrm{CuO}_3`$ for all-optical switching.
Nonlinear optical materials with large nonlinear coefficient, fast response-time, low loss and operability at room-temperature will be indispensable for next generation high speed network systems in which clock recovery and buffering at terabit/second rate will be performed by all-optical switches . These devices should operate in the transparent region below the fundamental band edge, where optical nonlinearity is mainly associated with two-photon absorption (TPA). Real carrier excitation, inevitable due to TPA, is considered as the major limiting factor here. Although terabit/second all-optical switches have been already demonstrated with conventional inorganic band semiconductors, considerable effort is necessary to overcome the difficulty associated with long carrier lifetime. Large optical nonlinearity and sub-picosecond response times are observed in organic $`\pi `$-conjugated polymers , but further improvements in sample quality and morphology would be required for actual applications. Parallel development of novel nonlinear optical inorganic materials that on the one hand possess the large nonlinearity and ultrafast recovery times of the organics, and the intrinsic robustness and superior thermal conductivity of inorganics on the other, is highly desirable. To date, no such inorganic semiconductor is known. We report in the present Letter the observation of large ultrafast optical nonlinearity in a novel strongly correlated inorganic semiconductor that is intrinsically different from conventional inorganic band semiconductors. The specific material we have studied is $`\mathrm{Sr}_2\mathrm{CuO}_3`$.
$`\mathrm{Sr}_2\mathrm{CuO}_3`$ is a quasi-1D Cu-O linear chain compound, whose structure is shown in Fig. 1(a). The material is a prototypical strongly correlated charge-transfer insulator within the Zaanen-Sawatzky-Allen scheme . Optical absorption in $`\mathrm{Sr}_2\mathrm{CuO}_3`$ involves charge-transfer (CT) excitation of a Cu-hole to an O-site, with an absorption band-edge at $``$ 1.6 eV and the band maximum at 2 eV (Fig. 1(b)). In addition to the high energy charge excitations, there exist low energy spin excitations in $`\mathrm{Sr}_2\mathrm{CuO}_3`$ , which correspond to those of a nearly ideal 1D Heisenberg chain with intrachain exchange integral $`J20003000\mathrm{K}`$. The interchain coupling in this system is rather weak, as is evidenced by the occurrence of 3D long-range antiferromagnetic coupling only below $`T_N=5\mathrm{K}`$ .
We present here the results of time-resolved femtosecond pump-probe measurements on a single crystal of $`\mathrm{Sr}_2\mathrm{CuO}_3`$, grown by the traveling-solvent floating-zone method . A thin flake with a thickness of $`L50\mu \mathrm{m}`$ was cleaved out with the $`bc`$-plane for transmission measurement. Optical pulses are provided by a system based on amplified mode-locked Ti:sapphire laser and optical parametric generator supplemented with sum and difference frequency generation. This system generates pulses with the photon energy centered at 0.2 – 2.0 eV and with temporal width $`200\mathrm{fs}`$. In our experiment, we measure the differential transmission $`\mathrm{\Delta }T/T`$ (where $`T`$ is the transmission in the absence of the pump beam), as a function of the delay time of the probe with respect to the pump.
Fig. 1(c) shows the temporal evolution of $`\mathrm{\Delta }T/T`$ measured at $`10K`$ for different pump ($`\omega _{\mathrm{pump}}`$) and probe ($`\omega _{\mathrm{probe}}`$) photon energies. The intensity of the pump pulse is $`2\mathrm{GW}/\mathrm{cm}^2`$. The intensity of the probe pulse is one to two orders of magnitude smaller. One can observe from Fig. 1(c) that the differential transmission consists of two components. The first component manifests itself as a sharp peak-like structure with temporal width of about $`200\mathrm{fs}`$, which is given by the overlapping of pump and probe pulses. This component, which is pronounced when the pump photon energy is tuned below the absorption edge of $`1.6\mathrm{eV}`$, is due to coherent interaction between the pump and probe pulses. At $`\omega _{\mathrm{pump}}`$= $`3.1\mathrm{eV}`$, which is far above the absorption edge, this coherent interaction is suppressed by the strong linear absorption, and correspondingly, the peak-like structure in the pump-induced transmission vanishes. At $`\omega _{\mathrm{pump}}`$ = $`1.55\mathrm{eV}`$, the strong group velocity dispersion near the absorption edge gives rise to the increase in the temporal width of the peak-like structure up to approximately $`1\mathrm{ps}`$.
The second component of the transmission change manifests itself as an exponentially decaying tail with characteristic time of about $`2\mathrm{ps}`$. One can observe from Fig. 1(c) that magnitude of this component decreases with decreasing of the pump photon energy and, correspondingly, with the pump absorption coefficient (see Fig. 1(b)). These experimental findings indicate that this relatively long-lived component of the transmission change can be associated with real excitation of the CT exciton, whose decay time is of similar magnitude.
These features in the temporal behavior of the pump-induced transmission remain at room temperature. We find that the magnitude of the effect is proportional to the pump intensity up to $`5\mathrm{GW}/\mathrm{cm}^2`$, indicating the pump-induced transmission change to be a third order optical process.
We describe $`\mathrm{\Delta }T/T`$ in terms of the pump-induced change in the absorption coefficient $`\mathrm{\Delta }\alpha =\mathrm{ln}(1+\mathrm{\Delta }T/T)/L`$, where $`L`$ is the crystal length, which consists of the instantaneous (peak) and exponentially decaying (tail) components, respectively: $`\mathrm{\Delta }\alpha =\mathrm{\Delta }\alpha _{\mathrm{peak}}+\mathrm{\Delta }\alpha _{\mathrm{tail}}`$. Fig. 2(a) shows the spectrum of TPA coefficient $`\beta =\mathrm{\Delta }\alpha _{\mathrm{peak}}/I_{\mathrm{pump}}`$ as a function of $`\omega _{\mathrm{probe}}`$ for pump photon energies of $`0.7\mathrm{eV}`$, $`1.1\mathrm{eV}`$ and $`1.55\mathrm{eV}`$ . One can observe from Fig. 2(a) that the smaller the $`\omega _{\mathrm{pump}}`$, the larger is the probe photon energy $`\omega _{\mathrm{probe}}`$ at which the maximum of the pump-induced coherent absorption is obtained. Moreover, it is seen that $`\beta `$ is maximum at $`\omega _{\mathrm{probe}}2.1\mathrm{eV}`$$`\omega _{\mathrm{pump}}`$ at all pump photon energies. This strongly suggests that the peak structure in Fig. 1 (c) is due to a two-photon allowed (one-photon forbidden) state at $``$ 2.1 eV. Alternative mechanisms, which may be responsible for the observed instantaneous photoinduced absorption change, can in principle involve fifth- and higher order effects associated with real carriers excitation due to two-photon absorption of the pump or with the electron-hole plasma created by the intense pump pulse . However, the observed linear dependence of $`\mathrm{\Delta }\alpha _{\mathrm{peak}}`$ on the pump intensity suggests that the instantaneous photoinduced absorption is essentially of the third-order in the light field. This observation along with the strong dependence of the pump-induced transmission spectra on $`\omega _{\mathrm{pump}}`$ precludes mechanisms of optical nonlinearity other than TPA.
To clarify this further we plot $`\beta `$ in Fig. 2(b) as a function of $`\omega _{\mathrm{pump}}+\omega _{\mathrm{probe}}`$. One can readily observe that maximum of the pump-induced coherent absorption takes place at $`\omega _{\mathrm{pump}}+\omega _{\mathrm{probe}}=2.1\mathrm{eV}`$ for all pump photon energies. The dashed line on Fig. 2(b) represents the linear absorption spectrum indicating that one- and two-photon allowed bands nearly overlap.
In order to clarify the origin of the observed TPA band in Fig. 2(b) we consider the Cu-O chain within the two-band extended Hubbard model ,
$$H=\underset{i}{}U_in_in_i+V\underset{i}{}n_in_{i+1}+\underset{i,\sigma }{}(1)^iϵn_it\underset{i,\sigma }{}(c_{i\sigma }^{}c_{i+1\sigma }+c_{i+1\sigma }^{}c_{i\sigma })$$
(1)
where $`c_{i\sigma }^{}`$ creates a hole with spin $`\sigma `$ on site $`i`$, $`n_{i\sigma }=c_{i\sigma }^{}c_{i\sigma }`$, and $`n_i=_\sigma n_{i\sigma }`$. The parameter $`t`$ is the transfer integral between $`Cu`$ and $`O`$ sites, $`2ϵ=ϵ_Oϵ_{Cu}`$, where $`ϵ_O`$ and $`ϵ_{Cu}`$ are the site energies of the $`O`$ and $`Cu`$ sites, $`U_i`$ is the on-site Coulomb repulsion between two holes on $`Cu`$ and $`O`$ sites ($`U_{Cu}U_O`$) and $`V`$ is the Coulomb repulsion between holes occupying neighboring $`Cu`$ and $`O`$ sites. For large $`U_i`$ and realistic $`ϵ>0`$, the holes occupy predominantly the $`Cu`$ sites in the ground state, which can thus be represented by the hole occupancy scheme ..101010…., where 1(0) represents an occupied(unoccupied) site. From the nature of the current operator $`\widehat{ȷ}=it_{i,\sigma }(c_{i\sigma }^{}c_{i+1\sigma }c_{i+1\sigma }^{}c_{i\sigma })`$, the one-photon optical state has the form (…10100110…) – (…10110010…), i.e., the odd parity linear combination of the two configurations that are reached by one application of the current operator on the ground state. The even parity “plus” linear combination of the same two configurations is a two-photon state, which for strong correlations is nearly degenerate with the one-photon state. This near degeneracy is expected to lead to very large transition dipole coupling between the one- and two-photon states.
To confirm the above conjecture we have performed exact numerical calculations for finite periodic rings of 12 sites (6 $`Cu`$ and 6 $`O`$) within Eq.1, for parameters $`|t|`$ = 1 - 1.4 eV, $`U_{Cu}`$ = 8 -10 eV, $`U_O`$ =4 - 6 eV, $`V`$ = 0 - 2 eV, and $`ϵ`$ = 1 – 2 eV. Because of the large Hubbard $`U`$ the relevant wavefunctions are strongly localized, and even the 12-site periodic ring can give semi-quantitative results. In all cases, the transition dipole coupling between the one- and the two-photon states is one to two orders of magnitude larger than that between the ground state and the one-photon state. The third order nonlinear optical coefficient, $`\chi ^{(3)}(\omega ;\omega ,\omega ,\omega )`$, whose imaginary component gives the TPA, is now calculated from the energies and transition dipole couplings. In Fig. 3 we have shown the calculated absorption and TPA spectra in arbitrary units, for one set of parameters for our finite system. The inset shows the energies of the one- and two-photon states and the transition dipole couplings that are used in the calculation.
At the maximum of the TPA band, the magnitude of the experimental TPA coefficient $`\beta `$ is $`150\mathrm{cm}/\mathrm{GW}`$ (Fig. 2(b)), which corresponds to $`\mathrm{Im}\chi ^{(3)}10^9`$ esu. This value is larger than that predicted from the gap-dependent scaling law derived for conventional semiconductors by one order of magnitude, and is comparable to that of conventional 1D-structured materials. The scaling law is inapplicable to Mott-insulators, in which the origin of the optical nonlinearity is the very large dipole coupling between nearly degenerate excited states of opposite parities. The mechanism of optical nonlinearity here is related to that in the $`\pi `$-conjugated polymers, which are described within the one-band extended Hubbard model and in which also there occurs very large dipole coupling between the optical state and a two-photon state slightly higher in energy . Not surprisingly, the magnitude of $`\chi ^{(3)}`$ in $`\mathrm{Sr}_2\mathrm{CuO}_3`$ is therefore comparable to some of the best organic materials . The intensity dependent refractive index $`n_2`$, obtained by a Kramers-Kronig transformation of the TPA data, is $``$ 10<sup>-7</sup> – 10<sup>-6</sup> cm<sup>2</sup>/MW, also comparable to the organics .
In addition to the large magnitude, the excitonic nonlinearity in $`\mathrm{Sr}_2\mathrm{CuO}_3`$ is featured by picosecond response, whose characteristic time is given by the decay constant of the tail-like component in Fig. 1(c). Such an ultrashort response indicates the existence of a fast non-radiative relaxation channel. Although detailed discussion of relaxation mechanisms of the optical excitation in $`\mathrm{Sr}_2\mathrm{CuO}_3`$ is beyond the scope of this paper, we suggest that the ultrafast ground state recovery is related to the occurrence of spin excitations below the optical gap. The low energy excitations of these system are the gapless spin excitations (spinons) of the uniform one-dimensional antiferromagnetic Heisenberg chain. These spin excitations are optically silent, and have an overall bandwidth of $``$ 1 eV . It is then conceivable that following the relaxation of the CT exciton to the highest energy spin excited states there occur further fast intra-spinon-band relaxation through the emission of multiple phonons and spinons. Because of the absence of such midgap states in the conventional semiconductors similar non-radiative processes would be absent there.
The room temperature ultrafast ground state recovery implies a high potential of $`\mathrm{Sr}_2\mathrm{CuO}_3`$ in the ultrafast optoelectronics and, specifically, in all-optical switching devices. In order to estimate this potential in terms of the maximum available repetition rate, we have performed a double-pulse experiment at room-temperature. Fig. 4 shows the $`\mathrm{\Delta }T/T`$ induced by two pump pulses with wavelength of $`1400\mathrm{nm}`$. The transmission change was measured at $`1200\mathrm{nm}`$, which is around the optical fiber communication wavelength. The second pump pulse, applied 2 ps after the first, induces nearly the same transmission change as the first. However, comparing with the single-pulse response (not shown), the transmission change in tail part is accumulated from pulse to pulse, leading to the limitation on the repetition rate. If $`\mathrm{\Delta }\alpha _{\mathrm{tail}}`$ and $`\mathrm{\Delta }\alpha _{\mathrm{peak}}`$ are the tail and peak components of photoinduced absorption and $`f`$ is the pulse repetition rate, the tail component of the induced absorption in the steady-state regime is $`\delta \alpha _{\mathrm{tail}}=\mathrm{\Delta }\alpha _{\mathrm{tail}}/[1\mathrm{exp}(1/\tau f)]`$. Therefore, the maximum available repetition rate can be estimated from a natural criterion of the operability of the system, $`\delta \alpha _{\mathrm{tail}}=\mathrm{\Delta }\alpha _{\mathrm{peak}}`$, which gives $`f_{\mathrm{max}}=(\mathrm{\Delta }\alpha _{\mathrm{peak}}/\mathrm{\Delta }\alpha _{\mathrm{tail}})\tau ^1`$. \>From Fig. 4, we obtain $`f_{\mathrm{max}}10^{13}\text{s}^1`$ for our sample, i.e., it can be used as a nonlinear optical medium with operability of several terabits per second.
In summary, we find that $`\mathrm{Sr}_2\mathrm{CuO}_3`$ exhibits strong nonlinearity and picosecond recovery of optical trasparency. Theoretical calculations indicates that the nonlinearity of the quasi-1D cuprates is due to a very large transitional dipole moment between nearly degenerate one- and two-photon states. Our findings suggest a strong potential of these materials for high bit-rate all-optical switching and, therefore, introduce a new means of achieving ultrafast optoelectronics with strongly correlated electron systems. We believe that with innovative material processing there is considerable scope for future enhancement of the figure of merit of these materilals for optoelectronic applications.
This work was supported by a grant-in-aid for COE Research from the Ministry of Education, Science, Sports, and Culture of Japan, and the U.S. NSF-ECS. |
warning/0002/hep-ph0002095.html | ar5iv | text | # Covariant Description of Composite Meson Systems and Chiral Symmetry
## 1 Introduction
There are the two contrasting view points of composite quark-antiquark mesons: The one is non-relativistic, based on the approximate symmetry of $`LS`$-coupling in NRQM; while the other is relativistic, based on the dynamically broken chiral symmetry in the NJL model. The $`\pi `$-meson (or $`\pi `$-nonet) is now widely believed to have a dual nature of non-relativistic particle with $`(L,S)=(0,0)`$ and also of relativistic particle as a Nambu-Goldston boson with $`J^P=0^{}`$ in the case of spontaneous breaking of chiral symmetry. However, no attempts to unify the above two view points have been yet proposed. On the other hand we have developed the covariant oscillator quark model (COQM) for many years as a covariant extension of NRQM, which is based on the boosted $`LS`$-coupling scheme. The meson wave functions (WF) in COQM are tensors in the $`\stackrel{~}{U}(4)O(3,1)`$ space and reduce at the rest frame to those in the $`SU(2)_{\mathrm{spin}}O(3)_{\mathrm{orbit}}`$ space in NRQM. The COQM has been successful especially in treating the $`Q\overline{Q}`$ meson system and the ($`q,\overline{Q}`$) or ($`Q,\overline{q}`$) meson system, leading, respectively, to a satisfactory understanding of radiative transitions and to the same weak form factor relations as in HQET. However, in COQM no consideration on chiral symmetry has been given and it is not able to explain the dual nature of $`\pi `$ meson.
The purpose of the present talk is to get rid of this defect in COQM and is to give a unified view point of the two contrasting view points of the composite meson systems.
## 2 Covariant Framework for Describing Composite Mesons
For meson WF described by $`\mathrm{\Phi }_A{}_{}{}^{B}(x_1,x_2)`$ ($`x_1,x_2`$ denoting the space-time coordinate and $`A=(\alpha ,a)(B=(\beta ,b))`$ denoting the Dirac spinor and flavor indices of constituent quark (anti-quark)) we set up the bilocal Yukawa equation
$`[^2/X_\mu ^2`$ $``$ $`^2(x_\mu ,/x_\mu )]\mathrm{\Phi }_A^B(X,x)=0`$ (1)
($`X(x)`$ denoting the center of mass (CM) (relative) coordinate of meson), where the $`^2`$ is squared mass operator including only a central, (Dirac) spinor
-independent confining potential. The WF is separated into the plane wave describing CM motion and the internal WF as
$`\mathrm{\Phi }_A{}_{}{}^{B}(x_1,x_2)={\displaystyle \underset{𝐏_n,n}{}}(e^{iP_nX}\mathrm{\Psi }_{n,A}{}_{}{}^{(+)B}(x,P_n)+e^{iP_nX}\mathrm{\Psi }_{n,A}{}_{}{}^{()B}(x,P_n)),`$ (2)
where $`P_{n,\mu }^2=M_n^2,P_{n,0}=\sqrt{M_n^2+𝐏_n^2}`$; $`(\pm )`$ represents the positive (negative) frequency part; and $`n`$ does a freedom of excitation. We have the following field theoretical expression in mind for the WF and its Pauli-conjugate, $`\overline{\mathrm{\Phi }}_B{}_{}{}^{A}(x_2,x_1)[\gamma _4\mathrm{\Phi }(x_1,x_2)^{}\gamma _4]_B^A`$ as
$`\mathrm{\Phi }_A{}_{}{}^{B}(x_1,x_2)`$ $`=`$ $`{\displaystyle \underset{n}{}}[0|\psi _A(x_1)\overline{\psi }^B(x_2)|M_n+M_n|\psi _A(x_1)\overline{\psi }^B(x_2)|0],`$
$`\overline{\mathrm{\Phi }}_B{}_{}{}^{A}(x_2,x_1)`$ $`=`$ $`{\displaystyle \underset{n}{}}[M_n|\psi _B(x_2)\overline{\psi }^A(x_1)|0+0|\psi _B(x_2)\overline{\psi }^A(x_1)|M_n],`$ (3)
where $`\psi _A(\overline{\psi }^B)`$ denotes the quark field (its Pauli-conjugate) and $`|M_n`$ does the composite meson state. The internal WF is, concerning the spin-dependence, expanded in terms of a complete set $`\{W^{(i)}\}`$ of free bi-Dirac spinors of quarks and anti-quarks; and the Fierz-component meson WF is expressed as
$`\mathrm{\Psi }_A^{(\pm )B}(x,P_n)={\displaystyle \underset{(i)}{}}W_\alpha ^{(i)(\pm )\beta }(P_n)M_a^{(i)(\pm )b}(x,P_n),M^{(i)(\pm )}(x,P_n)=ϵ^{(i)}\overline{W}^{(i)()}\mathrm{\Psi }^{(\pm )},`$ (4)
where $`A`$ means trace of $`A`$ and ortho-normal relations $`\overline{W}^{(i)()}W^{(j)(\pm )}=ϵ^{(i)}\stackrel{~}{\delta }_{ij}`$ ($`ϵ^{(i)}`$ and $`\stackrel{~}{\delta }_{ij}`$ denote,respectively, the sign and pseudo-Kronecker symbols) is used.
The meson WF satisfies, as is seen from Eq.(3), the self-conjugate condition, leading to the (conventional) crossing rule (or substituion law) for the (composite) meson WF, as
$`\overline{\mathrm{\Phi }}_A{}_{}{}^{B}(x_1,x_2)=\mathrm{\Phi }_A{}_{}{}^{B}(x_1,x_2),\overline{M}_b^{(i)(\pm )a}(x,P_n)=M_a^{(i)()b}(x,P_n)^{}=M_b^{(i)(\pm )a}(x,P_n).`$ (5)
## 3 Complete Set of Spin Wave Function and Heavy-quark Meson Systems
We set up the conventional “free” Dirac spinors with four-momentum of composite meson itself $`P=P_M`$, $`D_{q,\alpha }(P)(u_{q,\alpha }^{(r)}(P),v_{q,\alpha }^{(r)}(P)(r=1,2))`$ for quarks and $`\overline{D}_{\overline{q}}{}_{}{}^{\beta }(P)(\overline{v}_{\overline{q}}^{(s)\beta }(P),\overline{u}_{\overline{q}}^{(s)\beta }(P)(s=1,2))`$ for anti-quarks. It is to be noted that all four spinors are necessary for both quarks and anti-quarks inside of mesons. Then the complete set of bi-Dirac spinors is given by<sup>1</sup><sup>1</sup>1 In the following we give only the expressions of $`(+)`$-frequency parts, and consider only the ground states of composite system, disregarding the relative coordinates.
$`\{W^{(i)(+)}(P)\}`$ $`:`$ $`U(P)=u_q^{(r)}(p_1)\overline{v}_{\overline{q}}^{(s)}(p_2)|_{p_{i,\mu }=\kappa _iP_\mu }=u_+^{(r)}(𝐏)\overline{v}_+^{(s)}(𝐏),`$ (6)
$`C(P)=u_q^{(r)}(p_1)\overline{u}_{\overline{q}}^{(s)}(p_2)|_{p_{i,\mu }=\kappa _iP_\mu }=u_+^{(r)}(𝐏)\overline{v}_{}^{(s)}(𝐏),`$
$`D(P)=v_q^{(r)}(p_1)\overline{v}_{\overline{q}}^{(s)}(p_2)|_{p_{i,\mu }=\kappa _iP_\mu }=u_{}^{(r)}(𝐏)\overline{v}_+^{(s)}(𝐏),`$
$`V(P)=v_q^{(r)}(p_1)\overline{u}_{\overline{q}}^{(s)}(p_2)|_{p_{i,\mu }=\kappa _iP_\mu }=u_{}^{(r)}(𝐏)\overline{v}_{}^{(s)}(𝐏),`$
where $`u_+(𝐏)(\overline{v}_+(𝐏))`$ and $`u_{}(𝐏)(\overline{v}_{}(𝐏))`$ denote the Dirac spinors with positive energy and momentum $`𝐏`$ and with negative energy and momentum $`𝐏`$, respectively, describing quarks (anti-quarks). In Eq.(6) we have defined technically the momenta of constituent quarks (to be called quark-exciton) as
$`p_{i,\mu }\kappa _iP_\mu ,p_{i,\mu }^2=m_i^2;P_\mu ^2=M^2,M=m_1+m_2(\kappa _{1,2}m_{1,2}/(m_1+m_2)).`$ (7)
The respective members in Eq.(6) satisfy a couple of the corresponding free Dirac equations in momentum space (which are equivalent to the (conventional or new-type of) Bargman-Wigner Equations) and are expressed in terms of their irreducible composite meson WF as follows:
$`(\mathrm{Non}\mathrm{relat}.\mathrm{comp}.)`$ $`U_A{}_{}{}^{B}(P)={\displaystyle \frac{1}{2\sqrt{2}}}[(i\gamma _5P_{s,a}^{(NR)b}(P)+i\gamma _\mu V_{\mu ,a}^{(NR)b}(P))(1+{\displaystyle \frac{iP\gamma }{M}})]_\alpha {}_{}{}^{\beta },`$
$`(\mathrm{Semi}\mathrm{relat}.\mathrm{comp}.)`$ $`C_A{}_{}{}^{B}(P)={\displaystyle \frac{1}{2\sqrt{2}}}[(S_a^{(\overline{q})b}(P)+i\gamma _5\gamma _\mu A_{\mu ,a}^{(\overline{q})b}(P))(1{\displaystyle \frac{iP\gamma }{M}})]_\alpha {}_{}{}^{\beta },`$
$`D_A{}_{}{}^{B}(P)={\displaystyle \frac{1}{2\sqrt{2}}}[(S_a^{(q)b}(P)+i\gamma _5\gamma _\mu A_{\mu ,a}^{(q)b}(P))(1+{\displaystyle \frac{iP\gamma }{M}})]_\alpha {}_{}{}^{\beta },`$
$`(\mathrm{Extrly}.\mathrm{relat}.\mathrm{comp}.)`$ $`V_A{}_{}{}^{B}(P)={\displaystyle \frac{1}{2\sqrt{2}}}[(i\gamma _5P_{s,a}^{(ER)b}(P)+i\gamma _\mu V_{\mu ,a}^{(ER)b}(P))(1{\displaystyle \frac{iP\gamma }{M}})]_\alpha {}_{}{}^{\beta },`$ (8)
where all vector and axial-vector mesons satisfy the Lorentz conditions, $`P_\mu V_\mu (P)=P_\mu A_\mu (P)=0`$. Here it is to be noted that, in each type of the above members, the number of freedom counted both in the quark representation and in the meson representation is equal, as it should be ($`2\times 2=4`$ and $`1+3=4`$, respectively). Also it may be amusing to note that each constituent in all the above members is in “parton-like motion,” having the same 3-dimentional velocity as that of total mesons. (For example, in $`V(P),𝐯_{1,2}=\frac{𝐩^{(1,2)}}{p_0^{(1,2)}}=\frac{\kappa _{1,2}𝐏_M}{\kappa _{1,2}P_{M,0}}=𝐯_M`$.)
In the heavy quarkonium ($`Q\overline{Q}`$) system both quarks and anti-quarks are possible to do, since $`m_Q>\mathrm{\Lambda }_{\mathrm{conf}}`$, only non-relativistic motions with positive energy. Accordingly the bi-spinor $`U`$ is considered to be applied to $`Q\overline{Q}`$ system as a covariant spin WF. In the heavy-light quark meson $`Q\overline{q}`$($`q\overline{Q}`$) system the anti-quarks(quarks) make, since $`m_q\mathrm{\Lambda }_{\mathrm{conf}}`$, relativistic motions both with positive and negative energies. Accordingly both the bi-spinors $`U`$ and $`C`$ ($`U`$ and $`D`$) are to be applied to the $`Q\overline{q}`$($`q\overline{Q}`$) system, and in this system there is a possibility of existence of new composite scalar and axial-vector mesons(see Eq.(8)). In the light quark meson $`q\overline{q}`$-system both the quarks and anti-quarks make, since $`m_q\mathrm{\Lambda }_{\mathrm{conf}}`$, relativistic motions both with positive and negative energies. Accordingly the (linear combinations to be specified shortly of) bispinors $`U`$ and $`V`$ are applied to the $`q\overline{q}`$-system, and in this system there is a possibility of existence of an extra(, in addition, to a normal) set of composite pseudo-scalar and vector mesons.
## 4 Light-Quark Meson Systems and Chiral Symmetry
A)\[Charge conjugation\] properties of the bi-spinors and, correspondingly, of the composite mesons are derived from those of quarks as follows:
$`\mathrm{\Psi }_A^{(+)B}(P,x)`$ $`(`$ $`0|\psi _A(x_1)\overline{\psi }^B(x_2)|M)\mathrm{\Psi }_A^{c,(+)B}(P,x)(0|\psi _A(x_1)\overline{\psi }^B(x_2)|M^c)`$
$`=`$ $`(C^1)_{AA^{}}{}_{}{}^{t}\mathrm{\Psi }_{}^{(+)}(P,x)^A^{}{}_{B^{}}{}^{}C_{}^{B^{}B},(CC^{}=1,C\gamma _\mu C^1=^t\gamma _\mu );`$
$`P_{s,a}^{(NR)b}`$ $``$ $`P_{s,b}^{(NR)a},V_{\mu ,a}^{(NR)b}V_{\mu ,b}^{(NR)a};S_a^{(R\overline{q})b}S_b^{(Rq)a},A_{\mu ,a}^{(R\overline{q})b}A_{\mu ,b}^{(Rq)a};`$
$`P_{s,a}^{(ER)b}`$ $``$ $`P_{s,b}^{(ER)a},V_{\mu ,a}^{(ER)b}V_{\mu ,b}^{(ER)a}.`$ (9)
B)\[Chiral transformation\] properties of composite mesons are also derived straightforwardly from those of the bi-spinors
$`\psi _A{}_{}{}^{B}(P,x)`$ $``$ $`[e^{i\alpha ^i\frac{\lambda ^i}{2}\gamma _5}\psi (P,x)e^{i\alpha ^i\frac{\lambda ^i}{2}\gamma _5}]_A{}_{}{}^{B}.`$ (10)
C)\[Light-quark meson system-“chiral SU(6) multiplet”\]
The quark representation applying to the light-quark mesons is obtained by the linear transformation of the bi-Dirac spinors given in §3 as follows:
$`U_{P_s,\alpha }^{(N,E)\beta }`$ $``$ $`1/\sqrt{2}(U_{P_s}\pm V_{P_s})_\alpha {}_{}{}^{\beta }=i/2[\gamma _5(1,iv\gamma )]_\alpha {}_{}{}^{\beta };P_s^{(N,E)};C=(+,+)`$
$`C_{S,\alpha }^{(N,E)\beta }`$ $``$ $`(1,i)/\sqrt{2}(D_S\pm C_S)_\alpha {}_{}{}^{\beta }=1/2[(1,v\gamma )]_\alpha {}_{}{}^{\beta };S^{(N,E)};C=(+,)`$
$`U_{V,\alpha }^{(N,E)\beta }`$ $``$ $`1/\sqrt{2}(U_V\pm V_V)_\alpha {}_{}{}^{\beta }=i/2[(\stackrel{~}{\gamma }_\mu ,\sigma _{\mu \nu }v_\nu )]_\alpha {}_{}{}^{\beta };V^{(N,E)};C=(,)`$
$`C_{A,\alpha }^{(N,E)\beta }`$ $``$ $`(1,i)/\sqrt{2}(D_A\pm C_A)_\alpha {}_{}{}^{\beta }=i/2[\gamma _5(\stackrel{~}{\gamma }_\mu ,i\sigma _{\mu \nu }v_\nu )]_\alpha {}_{}{}^{\beta };A^{(N,E)};C=(+,)`$ (11)
($`v_\mu P_\mu /M,\stackrel{~}{\gamma }_\mu v_\mu 0`$; and $`U_{P_s}`$ denotes the coefficient bi-spinors of $`P_s`$ and so on ), where we have given also the charge-conjugation parity of the corresponding (hidden flavor) composite mesons. The chiral transformation properties of the new bi-spinors (for $`U_L(1)\times U_R(1)`$) are easily seen to be similar as the conventional ones as
$`1i\gamma _5,v\gamma \gamma _5v\gamma ,i\gamma _5\stackrel{~}{\gamma }_\mu i\stackrel{~}{\gamma }_\mu ,i\sigma _{\mu \nu }v_\nu \gamma _5\sigma _{\mu \nu }v_\nu .`$ (12)
D)\[“Local chiral SU(6) field”\]
Extending our considerations to include the $`()`$-frequency part, we are led to a unified expression of what to be called, Local Chiral SU(6) field, as
$`\mathrm{\Psi }_A{}_{}{}^{B}(X)`$ $`=`$ $`\mathrm{\Psi }_A^{(N)B}(X)+\mathrm{\Psi }_A^{(E)B}(X)`$
$`\mathrm{\Psi }_A^{(N)B}`$ $`=`$ $`1/2[i\gamma _{5\alpha }{}_{}{}^{\beta }P_{s,a}^{(N)b}+i\stackrel{~}{\gamma }_{\mu ,\alpha }{}_{}{}^{\beta }V_{\mu ,a}^{(N)b}+1_alpha{}_{}{}^{\beta }S_{a}^{(N)b}+(i\gamma _5\stackrel{~}{\gamma }_\mu )_\alpha {}_{}{}^{\beta }A_{\mu ,a}^{(N)b}]`$
$`\mathrm{\Psi }_A^{(E)B}`$ $`=`$ $`1/2[(i\gamma _5\gamma _\mu )_\alpha {}_{}{}^{\beta }M_{}^{1}_\mu P_{s,a}^{(E)b}+(\sigma _{\mu \nu })_\alpha {}_{}{}^{\beta }M_{}^{1}_\mu V_{\nu ,a}^{(E)b}`$ (13)
$`+i\gamma _{\mu ,\alpha }{}_{}{}^{\beta }M_{}^{1}_\mu S_a^{(E)b}+(i\gamma _5\sigma _{\mu \nu })_\alpha {}_{}{}^{\beta }M_{}^{1}_\mu A_{\nu ,a}^{(E)b}].`$
## 5 Concluding remarks
In this talk we have pointed out a possibility of existence of the new composite meson multiplets, when both or either constituent quark-excitons have “negative” energies and momenta. Here we should like to point out that the chiral symmetry of QCD guarantees this to be the case, and there exist some candidates (to be discussed elsewhere) for these new multiplets. |
warning/0002/cond-mat0002160.html | ar5iv | text | # Correlation Effects in Multi-Band Hubbard Model and Anomalous Properties of FeSi
## Acknowledgements
The authors would like to thank Professor H. Yamada for providing them the details of the band calculation (LMTO-ASA) for FeSi and for his useful comments. This work is supported by Grant-in-Aid for Scientific Research No.11640367 from the Ministry of Education, Science, Sports and Culture. |
warning/0002/hep-th0002221.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Chern-Simons gauge theory provides an excellent framework to study knot and link invariants. This framework has shown to be very useful in both, the study of polynomial invariants as the Jones polynomial and its generalizations, and the study of Vassiliev invariants or numerical invariants of finite type. Non-perturbative aspects of the theory play a fundamental role in the first context while perturbative ones are the main tool in the second. In this paper I will present a brief summary of the results on Vassiliev invariants achieved from perturbative Chern-Simons gauge theory in the last few years.
An important line of investigation in the context of Vassiliev invariants is the search of a universal combinatorial expression. Different approaches have been carried out. In the framework of perturbative Chern-Simons gauge theory explicit combinatorial formulae for all the primitive invariants up to order four have been obtained . The context based on this approach seems rather promising to obtain a general combinatorial expression.
As any other gauge theory, Chern-Simons gauge theory can be analyzed for different gauge fixings. Covariant gauges lead to complicated integrals expressions for Vassiliev invariants . Simpler integral expressions are obtained in non-covariant gauges of the light-cone type . However, for a non-covariant gauge fixing of the temporal type one obtains combinatorial expressions . It is in the last situation in the one that all intermediate integrals can be done so one ends with combinatorial expressions where only the information contained at the crossings is relevant.
Combinatorial expressions are much preferred to compute invariants and to study their properties. It turns out that the resulting combinatorial expressions can be rewritten in terms of Gauss diagrams. These facts notably simplify their explicit formulae. In this brief presentation I will collect all the expressions based on Gauss diagramas for all the primitive Vassiliev invariants up to order four.
An extended version of this presentation can be found in ref. . In the present paper, however, I include a table that was not presented in ref. because the corresponding computations were not carried out then. The appendix contains a table with all the primitive Vassiliev invariants up to order four for all prime knots with ten crossings, as computed from the combinatorial expressions provided by Chern-Simons gauge theory.
The paper is organized as follows. In sect. 2 I briefly describe the quantization procedure in the temporal gauge and comment on the essential ingredients which lead to the combinatorial expressions. In sect. 3 the resulting combinatorial expressions are presented in terms of Gauss diagrams. Finally, in sect. 4 aspects of future research directions are described.
## 2 Perturbative Chern-Simons gauge theory in the temporal gauge
I will begin reviewing the basic elements of Chern-Simons gauge theory. This theory is a quantum field theory whose action is based on the Chern-Simons form associated to a non-abelian gauge group. The fundamental data in Chern-Simons gauge theory are the following: a smooth three-manifold $`M`$ which will be taken to be compact, a gauge group $`G`$, which will be taken semi-simple and compact, and an integer parameter $`k`$. The action of the theory is the integral of the Chern-Simons form associated to a gauge connection $`A`$ corresponding to a gauge group $`G`$:
$$S_{\mathrm{CS}}(A)=\frac{k}{4\pi }_M\text{Tr}(AdA+\frac{2}{3}AAA).$$
(1)
The exponential of $`i`$ times this action is invariant under gauge transformations for integer $`k`$.
The metric-independence of the action (1) implies that the resulting quantum field theory is topological. Appropriate observables lead to vacuum expectation values (vevs) which correspond to topological invariants. A particularly important set of observables is constituted by Wilson loops. They correspond to the holonomy of the gauge connection $`A`$ along a loop. Given a representation $`R`$ of the gauge group $`G`$ and a 1-cycle $`\gamma `$ on $`M`$, it is defined as:
$$W_\gamma ^R(A)=\text{Tr}_R\left(\text{Hol}_\gamma (A)\right)=\text{Tr}_R\text{P}\mathrm{exp}_\gamma A.$$
(2)
Products of these operators are the natural candidates to obtain topological invariants after computing their vev. These vevs are formally written as:
$$W_{\gamma _1}^{R_1}W_{\gamma _2}^{R_2}\mathrm{}W_{\gamma _n}^{R_n}=[DA]W_{\gamma _1}^{R_1}(A)W_{\gamma _2}^{R_2}(A)\mathrm{}W_{\gamma _n}^{R_n}(A)\text{e}^{iS_{\mathrm{CS}}(A)},$$
where $`\gamma _1`$, $`\gamma _2`$,$`\mathrm{}`$,$`\gamma _n`$ are 1-cycles on $`M`$ and $`R_1`$, $`R_2`$ and $`R_n`$ are representations of $`G`$. In (2), the quantity $`[DA]`$ denotes the functional integral measure. The functional integral in (2) is not well defined. A variety of methods have been proposed to go around this problem and provide some meaning to the right hand side of (2). These methods fall into two categories, perturbative and non-perturbative ones. Witten, in his pioneer work 1988 , showed, using non-perturbative methods, that when one considers non-intersecting cycles $`\gamma _1`$, $`\gamma _2`$,$`\mathrm{}`$,$`\gamma _n`$ without self-intersections, the vevs (2) lead to polynomial invariants as the Jones polynomial and its generalizations.
The construction of the perturbative series expansion of the vev of an operator when dealing with a gauge theory starts with a gauge fixing. The first analysis of Chern-Simons perturbation theory were made in the covariant Landau gauge . Subsequent studies in this gauge led to a framework linked to the theory of Vassiliev invariants, which culminated with the configuration space integral approach .
Non-covariant gauges have been also studied in the context of Chern-Simons perturbation theory. The perturbative series which results in the non-covariant light-cone gauge leads to the Kontsevich integral . Vevs of gauge invariant operators are independent of the gauge chosen and therefore the expressions obtained in the light-cone gauge should be equivalent to the ones obtained in any other gauge. Thus, from a quantum field theory point of view, the configuration space integral which appear in the covariant Landau gauge leads to the same quantities as the Kontsevich integral. The non-covariant temporal gauge leads to alternative expressions for the vevs of gauge-invariant operators which turn out to be combinatorial . Again, gauge invariance implies that the resulting quantities must be the same as the ones in the configuration space approach or in the Kontsevich integral. In the rest of this section I will review the salient features of the analysis of the perturbative series expansion of the vacuum expectation value of a Wilson loop in the temporal gauge.
The gauge-fixing condition in the temporal gauge takes the form
$$n^\mu A_\mu =0,$$
(3)
where $`n`$ is the unit vector $`n^\mu =(1,0,0)`$. In this gauge the propagator becomes:
$$\delta _{ab}\frac{\lambda }{(np)^2}(p_\mu p_\nu \frac{i}{\lambda }(np)ϵ_{\mu \nu \rho }n^\rho )iϵ_{\mu \nu \rho }\frac{n^\rho }{np}\delta _{ab},$$
(4)
where the limit $`\lambda 0`$ has been taken. To construct the perturbative series expansion of the vev of a Wilson loop, one needs the Fourier transform of (4). In the temporal gauge the momentum-space integral that must be carried out has the form:
$$\mathrm{\Delta }(x_0,x_1,x_2)=_M\frac{\mathrm{d}^3p}{(2\pi )^3}\frac{\text{e}^{i(p^0x_0+p^1x_1+p^2x_2)}}{p^0}.$$
(5)
This integral is ill-defined due to the pole at $`p^0=0`$. To make sense of it a prescription has to be given to circumvent the pole. A precise prescription will not be taken. Instead, a rather general form will be used ,
$$\mathrm{\Delta }(x_0,x_1,x_2)=\frac{i}{2}\text{sign}(x_0)\delta (x_1)\delta (x_2)+f(x_1,x_2),$$
(6)
where $`f(x_1,x_2)`$ is a prescription-dependent distribution. The important consequence of the result (6) is that the dependence of $`\mathrm{\Delta }(x_0,x_1,x_2)`$ on $`x_0`$ has to be of the form $`\text{sign}(x_0)\delta (x_1)\delta (x_2)`$. This observation will be crucial in our analysis. The propagator (6) will allow us to introduce the notion of kernel of a Vassiliev invariant and to design a procedure to compute combinatorial expressions for these invariants.
Given a knot $`K`$ and one of its regular knot projections, $`𝒦`$, on the $`x_1,x_2`$-plane which is a Morse knot in the $`x_1`$ and $`x_2`$ directions, one has to deal with the following perturbative series expansion for the vacuum expectation value of the corresponding Wilson loop :
$$W(K,G)=W(𝒦,G)_{\mathrm{temp}}\times W(U,G)^{b(𝒦)},$$
(7)
being,
$$\frac{1}{d}W(K,G)=1+\underset{i=1}{\overset{\mathrm{}}{}}v_i(K)x^i,$$
(8)
and,
$$\frac{1}{d}W(𝒦,G)_{\mathrm{temp}}=1+\underset{i=1}{\overset{\mathrm{}}{}}\widehat{v}_i(𝒦)x^i.$$
(9)
In these expressions $`x`$ denotes the inverse of the Chern-Simons coupling constant, $`x=2\pi i/k`$, $`G`$ the gauge group, and $`d`$ the dimension of the representation carried by the Wilson loop. The function $`b(𝒦)`$ is the exponent of the Kontsevich factor, which has been conjectured to be ,
$$b(𝒦)=\frac{1}{12}(n_{x_1}+n_{x_2}),$$
(10)
where $`n_{x_1}`$ and $`n_{x_2}`$ are the critical points of the regular projection $`𝒦`$ in both, the $`x_1`$ and the $`x_2`$ directions. In (7) $`U`$ denotes the unknot and $`W(𝒦,G)_{\mathrm{temp}}`$ is the vacuum expectation of the Wilson loop corresponding to the regular projection $`𝒦`$ as computed perturbatively in the temporal gauge with the standard Feynman rules of the theory. Notice that though each of the factors on the right hand side of (7) depends on the regular projection chosen, the left hand side does not. While the coefficients $`v_i(K)`$ of the series (8) are Vassiliev invariants the coefficients $`\widehat{v}_i(K)`$ of (9) are not. The latter depend on the regular projection chosen.
## 3 Combinatorial expressions in terms of Gauss diagrams
A universal combinatorial formula for Vassiliev invariants could be obtained if the coefficients $`\widehat{v}_i(K)`$ in (9) could be computed with no integrals left. Unfortunately, this has not been obtained yet to all orders. Only part of the contributions entering $`\widehat{v}_i(K)`$ have been explicitly written to all orders. These are the kernels introduced in . The kernels are quantities which depend on the knot projection chosen and therefore are not knot invariants. However, at a given order $`i`$ a kernel differs from an invariant of type $`i`$ by terms that vanish in signed sums of order $`i`$. The kernel contains the part of a Vassiliev invariant which is the last in becoming zero when performing signed sums, in other words, a kernel vanishes in signed sums of order $`i+1`$ but does not in signed sums of order $`i`$. Kernels plus the structure of the perturbative series expansion seem to contain enough information to reconstruct the full Vassiliev invariants .
The expression for the kernels results after considering only the simplest part of the propagator of the gauge field in the temporal gauge. This part involves a double delta function and therefore all the integrals can be performed. The result is a combinatorial expression in terms of crossing signatures after distributing propagators among all the crossings. The general expression can be written in a universal form much in the spirit of the universal form of the Kontsevich integral . Let us consider a knot $`K`$ with a regular knot projection $`𝒦`$ containing $`n`$ crossings. Let us choose a base point on $`𝒦`$ and let us label the $`n`$ crossings by $`1,2,\mathrm{},n`$ as one passes for first time through each of them when traveling along $`𝒦`$ starting at the base point. The universal expression for the kernel associated to $`𝒦`$ has the form:
$`𝒩(𝒦)={\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}({\displaystyle \underset{m=1}{\overset{k}{}}}{\displaystyle \underset{\genfrac{}{}{0pt}{}{p_1,\mathrm{},p_m=1}{p_1+\mathrm{}+p_m=k}}{\overset{k}{}}}{\displaystyle \underset{\genfrac{}{}{0pt}{}{i_1,\mathrm{},i_m=1}{i_1<\mathrm{}<i_m}}{\overset{n}{}}}{\displaystyle \frac{ϵ_{i_1}^{p_1}\mathrm{}ϵ_{i_m}^{p_m}}{(p_1!\mathrm{}p_m!)^2}}`$
$`\text{ }{\displaystyle \underset{\genfrac{}{}{0pt}{}{\sigma _1,\mathrm{},\sigma _m}{\sigma _1P_1,\mathrm{},\sigma _mP_m}}{}}𝒯(i_1,\sigma _1;\mathrm{};i_m,\sigma _m)).`$
(11)
In this equation $`P_m`$ denotes the permutation group of $`p_m`$ elements. The factors in the innest sum, $`𝒯(i_1,\sigma _1;\mathrm{};i_m,\sigma _m)`$, are group factors which are computed in the following way: given a set of crossings, $`i_1,\mathrm{},i_m`$, and a set of permutations, $`\sigma _1,\mathrm{},\sigma _m`$, with $`\sigma _1P_1,\mathrm{},\sigma _mP_m`$, the corresponding group factor $`𝒯(i_1,\sigma _1;\mathrm{};i_m,\sigma _m)`$ is the result of taking a trace over the product of group generators which is obtained after assigning $`p_1,\mathrm{},p_m`$ group generators to the crossings $`i_1,\mathrm{},i_m`$ respectively, and placing each set of group generators in the order which results after traveling along the knot starting from the base point. The first time that one encounters a crossing $`i_j`$ a product of $`p_j`$ group generators is introduced; the second time the product is similar, but with the indices rearranged according to the permutation $`\sigma _jP_j`$.
The universal formula (11) for the kernels can be written in a more useful way collecting all the coefficients multiplying a given group factor. The group factors can be labeled by chord diagrams. At order $`k`$ one has a term for each of the inequivalent chord diagrams with $`k`$ chords. Denoting chord diagrams by $`D`$, equation (11) can be written as:
$$𝒩(𝒦)=\underset{D}{}N_D(𝒦)D,$$
(12)
where the sum extends to all inequivalent chord diagrams. The concept of kernel can be extended to include singular knots by considering signed sums of (12), or, following , introducing vacuum expectation values of the operators for singular knots. If $`𝒦^j`$ denotes a regular projection of a knot $`K^j`$ with $`j`$ simple singular crossings or double points, the corresponding universal form for the kernel possesses an expansion similar to (12):
$$𝒩(𝒦^j)=\underset{D}{}N_D(𝒦^j)D.$$
(13)
The general results about singular knots proved in lead to two important features for (13). On the one hand, finite type implies that $`N_D(𝒦^j)=0`$ for chord diagrams $`D`$ with more than $`j`$ chords. On the other hand, $`N_D(𝒦^j)=2^j\delta _{D,D(𝒦^j)}`$, where $`D(𝒦^j)`$ is the configuration corresponding to the singular knot projection $`𝒦^j`$. As observed above, kernels constitute the part of a Vassiliev invariant which survives a maximum number of signed sums.
To compute $`N_D(𝒦)`$ one needs to introduce first the notion of the set of labeled chord subdiagrams of a given chord diagram. This set will be denoted by $`S_D`$. This set is made out of a selected set of labeled chord diagrams that will be defined now. A labeled chord diagram of order $`p`$ is a chord diagram with $`p`$ chords and a set of positive integers $`k_1,k_2,\mathrm{},k_p`$, which will be called labels, such that each chord has one of these integers attached. The set $`S_D`$ is made out of labeled chord diagrams which satisfy two conditions. These conditions are fixed by the form of the series entering the kernels (11). The elements of $`S_D`$ will be called labeled chord subdiagrams of the chord diagram $`D`$. They are defined as follows. A labeled chord subdiagram of a chord diagram $`D`$ with $`k`$ chords is a labeled chord diagram of order $`p`$ with labels $`k_1,k_2,\mathrm{},k_p`$, $`pk`$, such that the following two conditions are satisfied: a) $`k_1+k_2+\mathrm{}+k_p=k`$; b) there exist elements $`\sigma _1P_{k_1},\sigma _2P_{k_2},\mathrm{},\sigma _pP_{k_p}`$ of the permutation groups $`P_{k_1},P_{k_2},\mathrm{},P_{k_p}`$ such that, after replacing the $`j`$-th chord diagram by $`k_j`$ chords arranged according to the permutation $`\sigma _j`$, for $`j=1,\mathrm{},p`$, the resulting chord diagram is homeomorphic to $`D`$. The number of ways that permutations $`\sigma _1P_{k_1},\sigma _2P_{k_2},\mathrm{},\sigma _pP_{k_p}`$ can be chosen is called the multiplicity of the labeled chord subdiagram. The multiplicity of a given labeled chord subdiagram, $`sS_D`$, will be denoted by $`m_D(s)`$.
The chord diagram $`D`$ itself can be regarded as a labeled chord subdiagram such that its labels, or positive integers attached to its chords, are 1. It has multiplicity 1. All the elements of $`S_D`$ except $`D`$ have a number of chords smaller than the number of chords of $`D`$. Not all labeled chord diagrams are subdiagrams of $`D`$. However, given a labeled chord diagram with labels $`k_1,k_2,\mathrm{},k_p`$ there can be different sets of permutations leading to $`D`$. The number of these different sets is the multiplicity introduced above. The elements of the sets $`S_D`$ for all chord diagrams $`D`$ up to order four which do not have disconnected subdiagrams are the following:
$``$ $`\text{ },`$
$``$ $`\text{ },\text{ },2`$
$``$ $`\text{ },\text{ },`$
$``$ $`\text{ },\text{ },\text{ },2`$
$``$ $`\text{ },2`$
$``$ $`\text{ },\text{ },2\text{ },4`$
$``$ $`\text{ },\text{ },2\text{ },`$
$``$ $`\text{ },\text{ },\text{ },2\text{ },2\text{ },3`$
$``$ $`\text{ },\text{ },\text{ },\text{ },`$
The numbers accompanying each labeled chord subdiagram denote their multiplicity. When no number is attached to a chord of a labeled chord diagram it should be understood that the corresponding label is 1.
In order to write the final expression for the kernels the notion of Gauss diagram must be introduced. Given a regular projection $`𝒦`$ of a knot $`K`$ we can associate to it its Gauss diagram $`G(𝒦)`$. The regular projection $`𝒦`$ can be regarded as a generic immersion of a circle into the plane enhanced by information on the crossings. The Gauss diagram $`G(𝒦)`$ consists of a circle together with the preimages of each crossing of the immersion connected by a chord. Each chord is equipped with the sign of the signature of the corresponding crossing. Gauss diagrams are useful because they allow to keep track of the sums involving the crossings which enter in (11) in a very simple form. Let us consider a chord diagram $`D`$ and one of its labeled chord subdiagrams $`sS_D`$. Let us assume that $`s`$ has $`p`$ chords and labels $`k_1,k_2,\mathrm{},k_p`$. We define the product,
$$s,G(𝒦),$$
(15)
as the sum over all the embeddings of $`s`$ into $`G(𝒦)`$, each one weighted by a factor,
$$\frac{\epsilon _1^{k_1}\epsilon _2^{k_2}\mathrm{}\epsilon _p^{k_p}}{(k_1!k_2!\mathrm{}k_p!)^2},$$
(16)
where $`\epsilon _1,\epsilon _2,\mathrm{},\epsilon _p`$ are the signatures of the chords of $`G(𝒦)`$ involved in the embedding. Using (15) the kernels $`N_D(𝒦)`$ entering (12) can be written as,
$$N_D(𝒦)=\underset{sS_D}{}m_D(s)s,G(𝒦),$$
(17)
where $`m_D(s)`$ denotes the multiplicity of the labeled subdiagram $`sS_D`$ relative to the chord diagram $`D`$.
The terms $`s,G(𝒦)`$ entering (17) are related to the quantities $`\chi (𝒦)`$ defined in . It is straightforward to obtain the following relations:
$`\text{ },G(𝒦)`$ $`=`$ $`{\displaystyle \frac{1}{(j!)^2}}\chi _1(𝒦),`$ $`\text{}\text{ },G(𝒦)`$ $`=`$ $`\chi _2^A(𝒦),`$ $`\text{}\text{ },G(𝒦)`$ $`=`$ $`{\displaystyle \frac{1}{16}}\chi _2^C(𝒦),`$ $`\text{}\text{ },G(𝒦)`$ $`=`$ $`\chi _3^B(𝒦),`$ $`\text{}\text{ },G(𝒦)`$ $`=`$ $`{\displaystyle \frac{1}{4}}\chi _3^D(𝒦),`$ $`\text{}\text{ },G(𝒦)`$ $`=`$ $`\chi _4^A(𝒦),`$ $`\text{}\text{ },G(𝒦)`$ $`=`$ $`\chi _4^C(𝒦),`$ $`\text{}\text{ },G(𝒦)`$ $`=`$ $`\chi _4^E(𝒦),\text{}`$ $`\text{ },G(𝒦)`$ $`=`$ $`{\displaystyle \frac{1}{(j!)^2}}n(𝒦),`$ $`\text{}\text{ },G(𝒦)`$ $`=`$ $`{\displaystyle \frac{1}{4}}\chi _2^B(𝒦),`$ $`\text{}\text{ },G(𝒦)`$ $`=`$ $`\chi _3^A(𝒦),`$ $`\text{}\text{ },G(𝒦)`$ $`=`$ $`{\displaystyle \frac{1}{4}}\chi _3^C(𝒦),`$ $`\text{}\text{ },G(𝒦)`$ $`=`$ $`{\displaystyle \frac{1}{4}}\chi _3^E(𝒦),`$ $`\text{}\text{ },G(𝒦)`$ $`=`$ $`\chi _4^B(𝒦),`$ $`\text{}\text{ },G(𝒦)`$ $`=`$ $`\chi _4^D(𝒦),`$ $`\text{}\text{ },G(𝒦)`$ $`=`$ $`\chi _4^F(𝒦),`$
where in the first row the relation in the left column applies when $`j`$ is odd and the one in the right when $`j`$ is even. Notice that in the second relation $`n(𝒦)`$ denotes the number of crossings of the regular projection $`𝒦`$.
In , Vassiliev invariants up to order four were expressed in terms of these quantities and the crossing signatures. The strategy to obtain them was to start with the kernels (17) and exploit the properties of the perturbative series expansion of Chern-Simons gauge theory. A special role in the construction was played by the factorization theorem proved in . Here, only their final form will be listed.
At second order, the final expression for the invariant is:
$$\alpha _{21}(K)=\alpha _{21}(U)+\text{ },G(𝒦)\text{ },G(\alpha (𝒦)),$$
(19)
where $`\alpha _{21}(U)`$ stands for the value of $`\alpha _{21}`$ for the unknot. In this expression $`\alpha (𝒦)`$ denotes the ascending diagram $`\alpha (𝒦)`$ of a knot projection $`𝒦`$. It is defined as the diagram obtained by switching, when traveling along the knot from a base point, all the undercrossings to overcrossings. Ascending diagrams enter often in the combinatorial expressions and it is convenient introduce the following notation. A bar over a quantity $`L(𝒦)`$ indicates that the same quantity for the ascending diagram has to be subtracted, i.e.:
$$\overline{L}(𝒦)=L(𝒦)L(\alpha (𝒦))$$
(20)
where $`\alpha (𝒦)`$ denotes the standard ascending diagram of $`𝒦`$. Using this notation, the final form for the only primitive Vassiliev invariant at order two is:
$$\alpha _{21}(K)=\alpha _{21}(U)+\text{ },\overline{G}(𝒦).$$
(21)
At order three there is only one primitive invariant. It takes the form:
$$\alpha _{31}(K)=\text{ }+\text{ }+2\text{ },G(𝒦)\underset{i=1}{\overset{n}{}}\epsilon _i(𝒦)\left[\text{ },G(\alpha (𝒦))\right]_i.$$
(22)
Several comments are in order to explain the quantities entering this expression. The sum is over all crossings $`i`$, $`i=1,\mathrm{},n`$, and $`\epsilon _i(𝒦)`$ denotes the corresponding signature. The square brackets $`[`$ $`]_i`$ enclosing a quantity $`L(𝒦)`$ denote:
$$\left[L(𝒦)\right]_i=L(𝒦)L(𝒦_{i_+})L(𝒦_i_{}),$$
(23)
where the regular projection diagrams $`𝒦_{i_+}`$ and $`𝒦_i_{}`$ are the ones which result after the splitting of $`𝒦`$ at the crossing point $`i`$.
Combinatorial expressions for the two primitive invariants at order four have been presented in . Their construction is based on the use of the kernels (17) and the properties of the perturbative series expansion. As in the case of previous orders, these invariants are expressed in terms of the products (15) and the crossing signatures. Their form is more complicated than the ones at lower orders. At order four there are two primitive Vassiliev invariants. The same choice of basis as in is made here. The combinatorial expressions for these two invariants turn out to be:
$`\alpha _{42}(K)=\alpha _{42}(U)+7\text{ }+5\text{ }+4\text{ }+2\text{ }+\text{ }+\text{ }`$
$`\text{ }+8\text{ }+2\text{ }+8\text{ }+{\displaystyle \frac{1}{6}}\text{ },\overline{G}(𝒦)`$
$`+{\displaystyle \underset{\genfrac{}{}{0pt}{}{i,j𝒞_a}{i>j}}{}}\overline{\epsilon }_{ij}(𝒦)(\left[\text{ },G(\alpha (𝒦))\right]_{ij}^a`$
$`\text{ }2\left[\text{ },G(\alpha (𝒦))\right]_i2\left[\text{ },G(\alpha (𝒦))\right]_j)`$
$`+{\displaystyle \underset{\genfrac{}{}{0pt}{}{i,j𝒞_b}{i>j}}{}}\overline{\epsilon }_{ij}(𝒦)(\left[\text{ },G(\alpha (𝒦))\right]_{ij}^b\left[\text{ },G(\alpha (𝒦))\right]_i`$
$`\text{ }\left[\text{ },G(\alpha (𝒦))\right]_j),`$
(24)
and,
$`\alpha _{43}(K)=\alpha _{43}(U)+\text{ }+\text{ }+\text{ }+2\text{ }{\displaystyle \frac{1}{6}}\text{ },\overline{G}(𝒦)`$
$`+{\displaystyle \underset{\genfrac{}{}{0pt}{}{i,j𝒞_a}{i>j}}{}}\overline{\epsilon }_{ij}(𝒦)(\left[\text{ },G(\alpha (𝒦))\right]_{ij}^a\left[\text{ },G(\alpha (𝒦))\right]_i`$
$`\text{ }\left[\text{ },G(\alpha (𝒦))\right]_j).`$
(25)
In these expressions the explicit dependence on the signatures appears in the quantities $`\overline{\epsilon }_{ij}(𝒦)`$ which are:
$$\overline{\epsilon }_{ij}(𝒦)=\epsilon _{ij}(𝒦)\epsilon _{ij}(\alpha (𝒦))=\epsilon _i(𝒦)\epsilon _j(𝒦)\epsilon _i(\alpha (𝒦))\epsilon _j(\alpha (𝒦)).$$
(26)
The sums in which these products are involved are over double splittings of the knot projection $`𝒦`$ at the crossings $`i`$ and $`j`$. There are two ways of carrying out these double splittings, depending on the configuration associated to the crossings $`i`$ and $`j`$. These are described in detail in . In the first one the regular projection $`𝒦`$ is split into two while in the second one it is split into three. Splittings of the first type build the set $`𝒞_a`$. The ones of the second type build $`𝒞_b`$. While only the first one is involved in the invariant $`\alpha _{43}`$, both appear in $`\alpha _{42}`$. The new quantities entering the sums are:
$`\left[L(𝒦)\right]_{ij}^a`$ $`=`$ $`L(𝒦)L(𝒦_{ij}^{a_1})L(𝒦_{ij}^{a_2}),`$
$`\left[L(𝒦)\right]_{ij}^b`$ $`=`$ $`L(𝒦)L(𝒦_{ij}^{b_1})L(𝒦_{ij}^{b_2})L(𝒦_{ij}^{b_3}),`$ (27)
where $`𝒦_{ij}^{a_1},𝒦_{ij}^{a_2},𝒦_{ij}^{b_1},𝒦_{ij}^{b_2}`$ and $`𝒦_{ij}^{b_3}`$ are the knot projections which originate after a double splitting of $`𝒦`$. As in previous orders, in the expressions (24) and (25), the quantities $`\alpha _{42}(U)`$ and $`\alpha _{43}(U)`$ correspond to the value of these invariants for the unknot. It has been proved in that the combinatorial expressions for $`\alpha _{42}(K)`$ and $`\alpha _{43}(K)`$ in (24) and (25) are invariant under Reidemeister moves.
Vassiliev invariants constitute vector spaces and their normalization can be chosen in such a way that they are integer-valued. Once their value for the unknot has been subtracted off they can be presented in many basis in which they are integers. We will chose here a particular basis in which the numerical values for the invariants up to order four are rather simple:
$`\nu _2(K)`$ $`=`$ $`{\displaystyle \frac{1}{4}}\stackrel{~}{\alpha }_{21}(K),`$
$`\nu _3(K)`$ $`=`$ $`{\displaystyle \frac{1}{8}}\stackrel{~}{\alpha }_{31}(K),`$
$`\nu _4^1(K)`$ $`=`$ $`{\displaystyle \frac{1}{8}}(\stackrel{~}{\alpha }_{42}(K)+\stackrel{~}{\alpha }_{43}(K)),`$
$`\nu _4^2(K)`$ $`=`$ $`{\displaystyle \frac{1}{4}}(\stackrel{~}{\alpha }_{42}(K)5\stackrel{~}{\alpha }_{43}(K)).`$ (28)
In these equations the tilde indicates that the value for the unknot has been subtracted, i.e., $`\stackrel{~}{\alpha }_{ij}(K)=\alpha _{ij}(K)\alpha _{ij}(U)`$. The values of the Vassiliev invariants (28) for all prime knots up to nine crossings have been presented in . The invariants (28) have been computed for torus knots in and . Denoting a generic torus knot by two coprime integers, $`p`$ and $`q`$, these invariants take the form:
$`\nu _2(p,q)`$ $`=`$ $`{\displaystyle \frac{1}{24}}(p^21)(q^21),`$
$`\nu _3(p,q)`$ $`=`$ $`{\displaystyle \frac{1}{144}}(p^21)(q^21)pq,`$
$`\nu _4^1(p,q)`$ $`=`$ $`{\displaystyle \frac{1}{288}}(p^21)(q^21)p^2q^2,`$ (29)
$`\nu _4^2(p,q)`$ $`=`$ $`{\displaystyle \frac{1}{720}}(p^21)(q^21)(2p^2q^23p^23q^23).`$
The explicit expression of Vassiliev invariants as polynomials in $`p`$ and $`q`$ is known up to order six . Of course, up to order four their value agree with the ones computed explicitly from equations (21), (22), (24) and (25). The only torus knots up to ten crossings are $`3_1`$, $`5_1`$, $`7_1`$, $`8_{19}`$, $`9_1`$ and $`10_{124}`$, whose associated coprime integers are (3,2), (5,2), (7,2), (4,3), (9,2) and (5,3), respectively.
In the table collected in the appendix the value of the primitive Vassiliev invariants for all the prime knots with ten crossings are presented. The value of these invariants for prime knots up to ten crossings can be found in .
## 4 Prospects
Though the perspectives are rather promising, the problems inherent to the proper treatment of gauge theories in non-covariant gauges do not permit at the moment to obtain a closed and complete formulation. Much work has to be done to understand the subtle issues related to the use of non-covariant gauges. The kernels plus the properties of the perturbative series expansion are probably enough to compute the explicit form of a given invariant but certainly it does not provide a systematic way of deriving the general universal formula. A proper and complete formulation of the perturbative series in a non-covariant gauge would explain the presence of the Kontsevich factor and will provide a general universal combinatorial formula. It is likely that a lattice formulation of Chern-Simons gauge theory in the temporal gauge could help considerably to make progress in this direction. We expect to report on this and other issues of perturbative Chern-Simons gauge theory in future work.
## Acknowledgments
I would like to thank the organizers of the Workshop on Geometry and Physics for inviting me to deliver a lecture and for their warm hospitality. This work is supported in part by DGICYT under grant PB96-0960.
## Appendix
In this appendix the result of computing the first four primitive Vassiliev invariants in (28), after using the expressions (21), (22), (24) and (25), for all prime knots with 10 crossings, is presented. These quantities have not been computed before using these combinatorial expressions. Their calculation is lengthy but straightforward once the computation of (21), (22), (24) and (25) are programmed and the prime knots are properly labeled. |
warning/0002/cond-mat0002427.html | ar5iv | text | # Shapes, contact angles, and line tensions of droplets on cylinders
## I Introduction
The wetting properties of a fiber in liquid matrices (e.g., dye mixtures, polymer melts, or molten resins) play an important role in the textile industry and in the fabrication of high-performance, fiber-reinforced composite materials. Since contact angles of liquid droplets on solid substrates provide a valuable characterization of such wetting properties there are numerous experimental and theoretical studies of the shape and the spreading of droplets deposited on a cylindrical substrate (see, e.g., Refs. ). The morphology of liquid drops on a fiber is particularly interesting insofar as on a planar substrate there is only one, spherical caplike droplet shape, whereas on a cylindrical substrate droplets may exhibit two, topologically different shapes, a “clamshell”- and a “barrel”-type one, depending on the droplet volume, the contact angle, and the cylinder radius . The aforementioned studies deal with thick fibers and large drops, i.e., the length scales are $`\mu `$m and larger. In this range the fluid structures are determined by macroscopic properties alone, i.e., volume of the liquid, surface tension $`\sigma `$ of the liquid vapor interface, Young’s contact angle $`\theta _{\mathrm{}}`$, and radius $`R`$ of the cylinder.
However, with the discovery of nanotubes the interest in such fluid structures has shifted to much smaller scales. There are several applications for which these small solid-fluid structures are very important. (i) For fabricating valuable composite materials involving nanotubes their wetting by the liquid host matrix is necessary to couple the inherent strength of the nanotubes to the matrix, reinforcing materials or fillers for plastics and ceramics . (ii) Nanotubes can be used as supports for heterogeneous catalysis or as templates for creating small wires or tubular structures by coating them with metals or metal oxides in the liquid state or by attaching inorganic and organic moieties to the nanotube surfaces . (iii) In order to use nanotubes as “nanostraws” potential candidates for exploiting such capillarity must be screened by first seeing if the liquid wets the *outside* of nanotubes . The performance of the nanotubes as catalysts, adsorbants, and deodorants can vary depending on whether they are composed of carbon, boron nitride, or oxides (SiO<sub>2</sub>, Al<sub>2</sub>O<sub>3</sub>, V<sub>2</sub>O<sub>5</sub>, MoO<sub>3</sub>, TiO<sub>2</sub>. This variety demonstrates, that the substrate potential of these tubes can be regarded as a tunable parameter. (iv) By using nanotubes as nanotweezers it might be possible to grab and manipulate small liquid drops. For this application the substrate must be nonwettable.
These small scales are comparable with the range of the substrate potential of the cylinders and of the molecular forces between the fluid particles adsorbing on them. Thus the droplets form under the action of the so-called effective interface potential $`\omega `$ which accounts for the net effect of the competition between the forces among the fluid particles and the substrate potential . Accordingly the calculation of the corresponding deformed droplet shapes requires a more detailed theoretical description which takes the effective interface potential into account. To our knowledge there is only one, recent publication in which this effect of $`\omega `$ on the droplet shape on fibers has been analyzed . It is the purpose of our study here to refine and to extend this analysis in various directions. If the radius $`R`$ of the fiber reduces to a few nm, as it is the case for nanotubes, the effective interface potential itself will depend on $`R`$ and thus deviate from that of the corresponding semi-infinite planar substrate used in Ref. . Accordingly we present a systematic analysis of the dependence of the shape of the droplets and their suitably defined contact angles on both $`R`$ and the droplet volume. This enables us to describe systematically the crossover in shape and contact angle between those of droplets on a cylinder and on the limiting case $`R\mathrm{}`$ of a planar substrate. We remark on how the structure of the effective interface potential, depending on whether it leads to first-order or continuous wetting transitions, influences the morphology of the droplets. We confine our analysis to barrel-type droplets and estimate their metastability against roll-up to the clamshell configuration. Finally we study two types of line tensions. The first one concerns the line tension of three-phase contact between liquid, vapor, and substrate emerging at the ends of macroscopicly large drops on fibers which reduces to the familiar line tension of the straight three-phase contact line on a planar substrate. The second excess free energy concerns the effective line tension associated with the circular shape of the three-phase contact line on a planar substrate as function of the droplet volume. These results are relevant for understanding how to extract line tensions from contact angle measurements.
We are encouraged to present our refined analyses by recent experimental advances to determine droplet shapes such as microscopic interferometry , ellipsometric microscopy , scanning polarization force microscopy , and tapping-mode scanning force microscopy . These techniques allow one to resolve drop profiles on the submicrometer scale down to the nanometer scale , both vertically and laterally. In view of the numerous important applications mentioned above it would be rather rewarding to extend the application of these techniques to nonplanar substrate geometries in order to resolve experimentally the shape of droplets on fibers and tubes as presented in the following sections.
## II Theory
### A Free energy functional
In cylindrical coordinates the droplet surface is described by a function $`h(z)`$ or $`l(z)`$ of the coordinate $`z`$ along the symmetry axis of the cylinder (Fig. 1(a)). We define $`h(z)`$ and $`l(z)`$ such that $`h(z)`$ is the local separation between the liquid-vapor interface and the symmetry axis of the cylinder and $`l(z)=h(z)R`$ is the local separation between the cylinder surface and the liquid-vapor interface, i.e., the liquid layer thickness. The droplet is also symmetric with respect to a reflection at the plane $`z=0`$. For large values of $`|z|`$, i.e., at large distances from the droplet center at $`z=0`$, the liquid forms a thin wetting layer of thickness $`l_0=h_0R`$ around the cylinder. For reasons of simplicity $`h(z)`$ is henceforth assumed to be a unique function of $`z`$, i.e., we do not consider contact angles $`\theta >90^{}`$. The shape of the liquid-vapor interface enclosing the droplet is determined by the interplay of three physical quantitites: the Laplace pressure $`2\sigma H`$ generated by the mean curvature $`H`$ of the interface with surface tension $`\sigma `$, the capillary pressure induced by the finite droplet volume, and the disjoining pressure or, equivalently, the effective interface potential $`\omega _c`$ acting on the liquid-vapor interface ; $`\omega _c(l;R)`$ is the cost in free energy per surface area to maintain a homogeneous wetting layer of prescribed thickness $`l`$ covering the cylinder surface and can be expressed in terms of the underlying forces of the substrate and between the fluid particles . In the absence of the effective interface potential, i.e., for large droplets the liquid-vapor interface is a minimal surface under the constraint of a prescribed volume, i.e., it exhibits a constant mean curvature. The influence of the effective interface potential is most pronounced near the cylinder surface within the range of the substrate potential and leads to a deviation of the actual profile $`h(z)`$ from the shape which is determined by the aforementioned constant mean curvature condition. On the other hand, in the limit of large separation from the cylinder surface the mean curvature is asymptotically constant because there the influence of the effective interface potential vanishes.
Independent of the size of the droplet, for later purposes we define the “reference configuration” (see Fig. 1(a))
$$a_{ref}(z)=a(z)\mathrm{\Theta }(z_1|z|)+h_0\mathrm{\Theta }(|z|z_1)$$
(1)
where $`a(z)`$ is that constant-mean-curvature surface which touches the surface $`h(z)`$, $`h(z=0)=a(z=0)`$, and exhibits the same curvature $`H`$ at the apex, i.e., the two principal radii of curvature
$$R_0=\frac{(1+(h^{}(0))^2)^{3/2}}{h^{\prime \prime }(0)}=\frac{1}{h^{\prime \prime }(0)}\text{with}h^{}\frac{dh}{dz}$$
(2)
and $`h(0)=R+l(0)`$ of $`h(z)`$ and, correspondingly, of $`a(z)`$ (see Fig. 1(a)) at the apex position $`z=0`$ are identical. $`\mathrm{\Theta }`$ denotes the Heaviside step function; $`\pm z_1`$ are those values of $`z`$ where $`a(z)`$ and the homogeneous wetting layer $`h(z)h_0`$ intersect. In this sense the values $`z=\pm z_1`$ define the positions of the two three-phase contact lines forming the ends of the droplet. The “apparent contact angle” $`\theta `$ is defined by the intersection of the barrel-shaped part $`a(z)`$ of the reference profile and the homogeneous wetting layer $`h(z)h_0`$ (see Fig. 1(a)):
$$\theta =\underset{zz_1}{lim}\mathrm{arctan}(|a^{}(z)|).$$
(3)
This apparent contact angle $`\theta `$ can be expressed in terms of the measurable quantities apex height $`l(0)`$ of the droplet, radius $`R`$ of the cylinder, radius $`R+l_0`$ of the wetting film, and radius of curvature $`R_0`$ (Eq. (2)) of the profile $`h(z)`$ at the apex:
$$\mathrm{cos}\theta =\frac{R+l(0)}{R+l_0}\frac{(R+l(0))^2(R+l_0)^2}{2(R+l_0)}\left(\frac{1}{R+l(0)}+\frac{1}{R_0}\right).$$
(4)
Within an interface displacement model (see, e.g., Ref. ) the equilibrium interface configuration $`\overline{h}(z)`$ for a droplet of prescribed *excess* volume $`V_{ex}`$ minimizes the free energy functional
$`F_{ex}[h(z)]`$ $`=`$ $`F[h(z)]F[h_0]`$ (5)
$`=`$ $`2\pi {\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}𝑑z\left(\sigma \left(h(z)\sqrt{1+h^{}(z)^2}h_0\right)+R\left(\omega _c(h(z)R)\omega _c(h_0R)\right)\right)`$ (6)
under the constraint
$$\pi \underset{\mathrm{}}{\overset{\mathrm{}}{}}𝑑z\left(h^2(z)h_0^2\right)=V_{ex}$$
(7)
and the boundary conditions $`h(|z|\mathrm{})=h_0`$. We have defined $`F_{ex}[h(z)]`$ as an excess free energy with respect to the free energy $`F[h_0]`$ of the homogeneous wetting layer $`h(z)h_0`$ rendering a mathematically well-defined, finite expression. The first contribution to $`F_{ex}`$ is the excess free energy due to the increase of the liquid-vapor interface as compared with a homogeneous cylindrical shape. In general $`\sigma `$ itself depends on the curvature and thus on $`R`$ (see, e.g., Sec. 2.2 in Ref. and references therein); in the following, however, we do not discuss explicitly this additional parametric dependence on $`R`$. The second contribution to $`F_{ex}`$ is the free energy generated by the effective interaction between the cylinder surface and the liquid-vapor interface, reduced by the corresponding free energy for the homogeneous wetting layer. Since the substrate is considered to be homogeneous, $`\omega _c(l)`$ depends only on the radial distance $`l=hR`$ from the substrate surface. The equilibrium separation $`h_0=l_0+R`$ of the homogeneous wetting layer from the cylinder axis minimizes the free energy $`F(h)=2\pi L(R\omega _c(hR)+\sigma h)`$ where $`Lz_1`$ is the macroscopic length of the cylinder. The constrained minimum of Eq. (5) is given by the unconstrained minimum of the surrogate functional
$$_{ex}[h(z)]=F_{ex}[h(z)]+\kappa \left(\pi \underset{\mathrm{}}{\overset{\mathrm{}}{}}𝑑z\left(h^2(z)h_0^2\right)V_{ex}\right).$$
(8)
The corresponding optimal profile $`\stackrel{~}{h}(z,\kappa )`$ renders the equilibrium profile $`\overline{h}(z,V_{ex})=\stackrel{~}{h}(z,\kappa (V_{ex}))`$ upon expressing the Lagrange multiplier $`\kappa `$ in terms of $`V_{ex}`$ by inserting $`\stackrel{~}{h}(z,\kappa )`$ into the left hand side of Eq. (7) which yields the implicit relation $`V_{ex}(\kappa )`$. In order to avoid a clumsy notation, in the following we denote $`\overline{h}(z,V_{ex})`$ by $`h(z)`$. The Euler-Lagrange equation corresponding to Eq. (8) reads
$$\sigma \left(\frac{1}{h(z)(1+h^{}(z)^2)^{1/2}}\frac{h^{\prime \prime }(z)}{(1+h^{}(z)^2)^{3/2}}\right)=2\sigma H(z)=\kappa \frac{R}{h(z)}\frac{d\omega _c(hR)}{dh}|_{h=h(z)}.$$
(9)
This equation describes the balance between the Laplace pressure on the left hand side and the capillary plus disjoining pressure on the right hand side. One has $`\kappa <0`$ for any barrel-shaped droplet.
### B Reference profiles
The reference profile $`a(z)`$ minimizes a similar surrogate functional:
$$𝒜[a(z)]=\pi \underset{z_1}{\overset{z_1}{}}𝑑z\left(2\sigma a(z)\sqrt{1+a^{}(z)^2}+\kappa ^{}a^2(z)\right)+\text{const}$$
(10)
with the constant independent of $`a(z)`$ and the boundary conditions $`a(\pm z_1)=h_0`$. Equation (10) follows from Eq. (8) by omitting $`\omega _c`$ and replacing $`\kappa `$ by $`\kappa ^{}`$. The corresponding Euler-Lagrange equation describes the constant-mean-curvature surface given by
$$\sigma \left(\frac{1}{a(z)(1+a^{}(z)^2)^{1/2}}\frac{a^{\prime \prime }(z)}{(1+a^{}(z)^2)^{3/2}}\right)=2\sigma H(z)=\kappa ^{}.$$
(11)
According to the definition of $`a(z)`$ the Lagrange multiplier $`\kappa ^{}`$ has to be chosen such that this constant mean curvature of this surface equals the mean curvature at the apex of the actual surface $`h(z)=l(z)+R`$:
$$\frac{\kappa ^{}}{\sigma }=\frac{1}{R+l(0)}+\frac{1}{R_0}$$
(12)
$`R_0`$ (see Eq. (2)) and $`R+l(0)`$ are the principal radii of curvature at the apex of the actual surface $`h(z)`$, determined by the former Lagrange multiplier $`\kappa (V_{ex})`$. The solution of Eq. (11) is given implicitly by
$$\underset{h_0}{\overset{a(z)}{}}𝑑y\left(\left(\frac{\sigma }{\kappa ^{}}\right)^2\left(\frac{2y}{𝒞y^2}\right)^21\right)^{1/2}=z+z_1,|z|z_1,$$
(13)
which fulfils the boundary condition $`a(z_1)=h_0`$; this determines implicitly $`z_1`$ in terms of $`h_0`$, $`V_{ex}`$, $`\sigma `$, and $`\omega _c`$. The integration constant $`𝒞`$ is determined by $`a^{}(0)=0`$ due to the symmetry of $`a(z)`$. The integral in Eq. (13) can be expressed in terms of elliptic integrals .
When the drop is macroscopicly large ($`V_{ex}=\mathrm{}`$) it is appropriate to adopt a slightly different point of view. In this case not the center of the droplet but the position of one of the two three-phase contact lines, which are defined by the intersection of the asymptote $`a_m(z)`$ (the constant-mean-curvature surface appertaining to the *m*acroscopic drop) and $`h(z)h_0`$, is fixed at $`z=0`$ (see Fig. 1(b)). The actual interface profile interpolates between, e.g., $`h(z\mathrm{})=h_0`$ and $`h(z\mathrm{})=a_m(z)`$. This configuration describes a solid cylinder which is in contact with bulk vapor on the left hand side ($`z\mathrm{}`$) and with bulk liquid on the right hand side ($`z\mathrm{}`$). The analysis of the internal structure of a three-phase contact line on a homogeneous, planar substrate is based on a similar configuration (see Refs. and and references therein). The interface profile diverges in the limit $`z\mathrm{}`$: $`h(z\mathrm{})\mathrm{}`$, but this divergence is not linear as in the case of the planar substrate. A macroscopicly large drop implies $`\kappa 0`$. In this limit the volume constraint loses its meaning. Instead the state of the system is fixed by different lateral boundary conditions. In this case the solution of Eq. (13) is given by
$$a_m(z)=a(z;\kappa =0)=C\text{cosh}\left(\frac{zD}{C}\right)$$
(14)
with two integration constants $`C`$ and $`D`$. $`a_m(z)`$ describes a rotational surface with minimal surface area. From Eq. (14) one can easily see that the divergence of the interface profile for macroscopic drops is *exponential*, $`a_m(z\mathrm{})=\frac{C}{2}\mathrm{exp}((zD)/C)`$, rather than linear as on a planar substrate. The reference profile appertaining to the macroscopic drop is
$$a_{ref,m}(z)=a_m(z)\mathrm{\Theta }(z)+h_0\mathrm{\Theta }(z).$$
(15)
The slopes at the intersection of the asymptote $`a_m(z)`$ and the homogeneous layer $`h(z)h_0`$ at $`z=0`$ defines the contact angle $`\theta _m(R)=\theta (R,V_{ex}\mathrm{})`$. $`R=\mathrm{}`$ corresponds to a planar substrate for which the interface profile diverges *linearly* in the limit $`z\mathrm{}`$:
$$a_{m,\mathrm{}}(z)R=l_0+z\mathrm{tan}\theta _{\mathrm{}}$$
(16)
with the macroscopic contact angle $`\theta _{\mathrm{}}=\theta _m(R\mathrm{})`$ on the planar substrate. $`\theta _{\mathrm{}}`$ obeys Young’s law $`\mathrm{cos}\theta _{\mathrm{}}=(\sigma _{wg}\sigma _{wl})/\sigma `$ where $`\sigma _{wg}`$ and $`\sigma _{wl}`$ are the wall-gas and wall-liquid surface tensions, respectively; $`\sigma _{wg}\sigma _{wl}=\omega _c(l_0;R=\mathrm{})`$ is determined by the effective interface potential of the corresponding planar substrate (see, c.f., Subsec. II C). On the cylindrical surface the contact angle $`\theta _m(R)`$ does not follow from similar thermodynamic considerations but follows from the numerical analysis of the full profile $`h(z)`$ for large $`V_{ex}`$ (see, c.f., Sec. III and Fig. 9).
The integration constants $`C`$ and $`D`$ in Eq. (14) can be determined from the conditions $`a_m(z=0)=R+l_0`$ and $`a^{}(z=0)=\mathrm{tan}\theta _m`$ so that
$$a_m(z)=R\mathrm{cos}\theta _m\mathrm{cosh}\left(\frac{z}{R\mathrm{cos}\theta _m}+\text{arccosh}\frac{1}{\mathrm{cos}\theta _m}\right).$$
(17)
The series expansion of this expression in terms of small $`z/R`$ is
$$a_m(z/R1)=R+l_0+z\mathrm{tan}\theta _m(R)+𝒪(z^2).$$
(18)
In the limit $`R\mathrm{}`$ the region where the higher order terms are relevant is shifted towards $`z=\mathrm{}`$ such that, with $`\theta _m(R\mathrm{})=\theta _{\mathrm{}}`$, one recovers the linearly diverging asymptote $`a_{m,\mathrm{}}(z)`$ (Eq. (16)).
### C Effective interface potential
For the same liquid layer thickness $`l`$ the effective interface potential $`\omega _c(l;R)`$ of a cylinder differs from that of a planar substrate $`\omega _p(l)`$. The full expression $`\omega _c(l;R)`$ is presented in Ref. as obtained from density functional theory and within a so-called sharp-kink approximation for the solid-liquid and the liquid-vapor interface profiles. For reasons of simplicity, here we use the leading order of a series expansion of $`\omega _c(l;R)`$ in terms of $`d_w/R`$ where $`d_w`$ is the radial extension of the volume excluded for the fluid particles due to the repulsive part of the substrate potential:
$`\omega _c(l;R)`$ $`=`$ $`{\displaystyle \frac{3\pi }{2}}a{\displaystyle \frac{R}{h^3}}_2F_1({\displaystyle \frac{5}{2}},{\displaystyle \frac{3}{2}};2;\left({\displaystyle \frac{R}{h}}\right)^2)+8b{\displaystyle \frac{R}{h^4}}_2F_1(3,2;2;\left({\displaystyle \frac{R}{h}}\right)^2)`$ (20)
$`+{\displaystyle \frac{315\pi }{32}}c{\displaystyle \frac{R}{h^9}}_2F_1({\displaystyle \frac{11}{2}},{\displaystyle \frac{9}{2}};2;\left({\displaystyle \frac{R}{h}}\right)^2)+𝒪\left({\displaystyle \frac{d_w}{R}}\right)`$
with $`h=l+R`$ and $`{}_{2}{}^{}F_{1}^{}`$ hypergeometric functions. In the limit $`l/R0`$ one recovers the expression for the effective interface potential of the corresponding planar substrate:
$$\omega _c(l;R\mathrm{})=\omega _c(l0,R)=\omega _p(l)=al^2+bl^3+cl^8.$$
(21)
However, the power-law decay of $`\omega _c(l\mathrm{})`$ for a fixed, finite cylinder radius $`R`$ is
$$\omega _c(l\mathrm{};R)=\frac{3\pi }{2}a\frac{R}{l^3}+𝒪(l^4),$$
(22)
i.e., one power faster than that for the corresponding planar substrate.
At present there exists, to our knowledge, only one study concerned with the shapes of droplets on cylinders within the range of the effective interface potential between the cylinder surface and the liquid-vapor interface . However, in Ref. the *disjoining pressure* $`\mathrm{\Pi }_c(l)=(R/(R+l))d\omega _c(l)/dl`$ on the right hand side of the Euler-Lagrange equation (9) as a whole rather than only the effective interface potential $`\omega _c`$ is replaced by the disjoining pressure of the corresponding planar substrate $`\mathrm{\Pi }_p(l)=d\omega _p(l)/dl`$. In view of Eqs. (21) and (22), except for the factor $`R/(R+l)`$, this corresponds to the short-distance expansion ($`l/R0`$) of the effective interface potential of the cylinder. This replacement of the disjoining pressure by that of the planar substrate is expected to yield numerically reliable results only for large cylinder radii and small liquid layer thicknesses $`lR`$. Therefore in Sec. III we test the quality of this approximation (as well as that of the replacement of $`\omega _c`$ alone by $`\omega _p`$).
So far, due to the volume constraint, our considerations apply to nonvolatile liquids. For volatile liquids any droplet surrounded by a macroscopic reservoir of vapor phase is thermodynamically unstable against evaporation, leaving behind only the thin equilibrium wetting film. However, we expect that the actual nonequilibrium state of a condensating or evaporating liquid observed within a time scale that is small compared with the typical condensation or evaporation time can be described by solutions of Eq. (9) with $`V_{ex}`$ given by its momentary value. Only the interface configuration for $`\kappa =0`$, i.e., $`V_{ex}=\mathrm{}`$, which interpolates between a homogeneous wetting layer and an exponentially diverging profile, describes a bona fide thermodynamically stable state which can be maintained by imposing appropriate boundary conditions (see above) at liquid-vapor coexistence for the bulk fluid. The thermodynamic state, which in a grand canonical ensemble is defined by temperature and chemical potential, enters parametrically into the actual values of the effective interface potential $`\omega _c`$ and the liquid-vapor surface tension $`\sigma `$.
## III Shapes of droplet surfaces and contact angles
We solve the Euler-Lagrange equation (9) numerically for fixed values of $`\kappa `$ and for a given effective interface potential $`\omega _c(l)`$; the value of $`\kappa `$, in turn, determines the excess liquid volume $`V_{ex}`$ and allows us to establish the relation $`\kappa (V_{ex})`$. As boundary conditions in the case $`\kappa <0`$ (leading to droplets of finite size) we use that $`h(z)`$ must approach the wetting layer thickness $`h_0`$ for large $`z`$ and that $`h^{}(z=0)=0`$. The distance $`L/2`$, at which the system is cut off, is chosen large enough so that $`h(z=L/2)`$ and $`h^{}(z=L/2)`$ attain their asymptotic values $`h_0`$ and $`0`$, respectively, within prescribed accuracy. The reference profile $`a_{ref}(z)`$ is then calculated numerically by solving the differential equation (11) with $`\kappa ^{}`$ determined by Eqs. (2) and (11) and with $`a(z=0)=h(z=0)`$ and $`a^{}(z=0)=0`$, up to the point of intersection of $`a(z)`$ and $`h_0`$ which defines the coordinate $`z_1`$; $`a_{ref}(zz_1)=h_0`$. The contact angle $`\theta `$ is determined from Eq. (3) and, as a crosscheck, from Eq. (4).
In all numerical calculations presented henceforth we set $`a=3\sigma s^2`$, $`b=5\sigma s^3`$, and $`c=3\sigma s^8`$ such that $`s`$ sets the length scale for the range of the effective interface potential (typically $`s1`$nm). We divide both sides of Eq. (9) by $`\sigma `$ so that $`\omega (l)/\sigma `$ is dimensionless and $`\kappa /\sigma `$ has the dimension of an inverse length. Alternatively, instead of introducing $`s`$ as above one can choose $`\sqrt{a/\sigma }`$ as the basic length scale which describes the decay of the effective interface potential; for our choice $`\sqrt{a/\sigma }1.73s`$. The effective interface potentials $`\omega _c(l;R)`$ (Eq. (20)) and $`\omega _p(l)`$ (Eq. (21)) for the above choice of coefficients are shown in Fig. 2.
As a first example we solve Eq. (9) with the effective interface potential $`\omega _p(l)`$ of the corresponding planar substrate (Eq. (21)) and the potential coefficients given above. Figure 3 shows the profile of the droplet surface on a cylinder with radius $`R=100s`$ for $`\kappa s/\sigma =0.1`$. This choice of $`\kappa `$ leads to a small droplet with $`V_{ex}1.46\times 10^4s^3`$ (i.e., containing roughly $`10^7`$ fluid particles) whose liquid-vapor interface lies entirely within the range of the effective interface potential. Therefore the deviation of the profile from the asymptote $`a(z)`$ extends up to the apex of the droplet. The model effective interface potential used here resembles a typical interface potential leading to first-order wetting on a planar substrate . The droplet surface crosses the reference profile and, upon approaching the apex of the droplet, it reaches the reference profile from below.
Carroll has shown that, in the absence of the effective interface potential, the axisymmetric droplet configuration is only stable for
$$2\left(\frac{h(0)}{R}\right)^3\mathrm{cos}\theta 3\left(\frac{h(0)}{R}\right)^2+1>0,$$
(23)
i.e., if the droplet is large compared with the diameter of the cylinder and if the contact angle is small. When the droplet volume decreases or the contact angle increases the axisymmetric droplet becomes metastable against a so-called “roll up” towards the “clamshell” configuration . Applying the stability criterion Eq. (23) to the interface profile shown in Fig. 3, we find that this barrel-type configuration is possibly metastable towards forming the “clamshell” shape. A definitive statement about the stability would require to refine the criterion in Eq. (23) by incorporating the effect of the effective interface potential. However, the determination of the non-axisymmetric “clamshell” equilibrium shape requires a much larger numerical effort and is therefore beyond the scope of the present paper. One can define a critical value $`V_{ex,c}`$ such that for $`V_{ex}>V_{ex,c}`$ the axisymmetric droplet is stable. Upon increasing $`R`$, $`V_{ex,c}`$ increases, too; $`V_{ex,c}\mathrm{}`$ in the limit $`R\mathrm{}`$. Only for $`V_{ex}=\mathrm{}`$, i.e., for macroscopic drops, and for contact angles smaller than $`90^{}`$ (as stated in Sec. II here we do not consider the case $`\theta >90^{}`$) the rotationally symmetric interface shape is stable for any value of $`R`$.
Figure 4 shows the droplet shape for the same choice of potential parameters and for the same cylinder with $`R=100s`$, but for $`\kappa s/\sigma =0.005`$. This choice of $`\kappa `$ leads to a much bigger droplet with $`V_{ex}1.67\times 10^8s^3`$. The apex of the droplet is located at such a large distance from the cylinder surface that the effective interface potential is almost negligible. Therefore the application of Eq. (23) is reliable; it yields that this particular droplet is indeed stable against “roll-up”. In the vicinity of the cylinder surface the absolute deviation of the interface profile from the asymptote is similar to that in Fig. 3. As compared with the situation shown in Fig. 3, the point where $`l(z)`$ crosses the reference profile $`a(z)R`$ is shifted to the right and lies near the three-phase contact line at $`z=z_1`$. For the model effective interface potential used here and in Fig. 3, in the region around the apex of the droplet the profile lies below the reference profile. These results are in accordance with the findings for the planar, homogeneous substrate with the same type of interface potential. If, on the other hand, $`\omega (l)`$ corresponds to a system undergoing a continuous wetting transition, i.e., exhibiting a single minimum without a potential barrier, the profile of the droplet shape approaches its asymptote from the outside without crossing it (see Fig. 8 in Ref. ).
Figure 5 displays the effect of the replacement of the effective interface potential of a cylinder $`\omega _c(l;R)`$ (Eq. (20)) by that of the corresponding planar substrate $`\omega _p(l)`$ (Eq. (21)) for a droplet with $`\kappa s/\sigma =0.005`$ so that $`l(0)361s`$ and $`z_1453s`$ on a thin cylinder with $`R=20s`$. This droplet also satisfies the stability criterion in Eq. (23). For reasons of clarity in this figure we have plotted the difference $`\mathrm{\Delta }h(z)`$ between corresponding profiles instead of the profiles themselves. The influence of approximating $`\omega _c(l;R)`$ by $`\omega _p(l)`$ turns out to be rather small: the difference $`\mathrm{\Delta }h(z)`$ is at most of the order of $`s`$. It is even smaller in the case of smaller droplets and thicker cylinders: the quality of approximating $`\omega _c(l;R)`$ by $`\omega _p(l)`$ improves if the cylinder is thicker and the droplet is smaller. In Ref. the disjoining pressure of a cylinder rather than the effective interface potential is replaced by its planar-substrate counterpart. This corresponds to replacing the term $`(R/h)d\omega _c/dl`$ on the right hand side of Eq. (9) by $`d\omega _p/dl`$. The dashed line in Fig. 5 shows the effect of this approximation on the surface profile. As expected, the quality of this approximation is worse than the substitution of the effective interface potential alone, although the difference $`\mathrm{\Delta }h`$ is still of the order of $`s`$.
Figure 6 shows the apparent contact angles for the examples presented in Figs. 3 and 4, as well as for the same set of potential parameters but with $`R=20s`$, $`R=200s`$, and $`R=500s`$ as function of the excess liquid volume $`V_{ex}`$. Upon increasing $`R`$ the curves are shifted upwards and to the right. Due to the exponential divergence of the interface profile of a macroscopic drop (which is more pronounced for smaller $`R`$, see the discussion of Eqs. (17) and (18) above and, c.f., Fig. 8), the determination of the apparent contact angles for very large drops is, in particular for thin cylinders, numerically difficult. However, the data indicate that, for any value of $`R`$, $`\theta (V_{ex})`$ is a monotonously decreasing function, with a vanishing slope $`d\theta /dV_{ex}=0`$ at $`V_{ex}=\mathrm{}`$. The differences between the contact angles for different $`R`$ are minimal at $`V_{ex}=\mathrm{}`$. For any $`R`$ there is a sizeable increase of the apparent contact angle upon decreasing droplet size.
Figure 7 shows the effect of replacing $`\omega _p`$ by $`\omega _c`$ on the apparent contact angles for the system with $`R=20s`$ (compare Fig. 5). The difference between the contact angles calculated by using $`\omega _c`$ and $`\omega _p`$ is significant. It is much smaller for thicker cylinders whose contact angles are displayed in Fig. 6. However, the qualitative functional form of the dependence of $`\theta (V_{ex})`$ is not affected by the replacement of $`\omega _c`$ by $`\omega _p`$.
For macroscopic drops (i.e., $`\kappa =0`$ or, equivalently, $`V_{ex}=\mathrm{}`$) the Euler-Lagrange equation is solved with the initial value $`h(z=L_1)=h_0`$ and a small initial slope $`h^{}(z=L_1)`$ (e.g., $`h^{}(z=L_1)=10^8`$); the initial slope $`h^{}(z=L_1)=0`$ would yield the trivial solution $`h=h_0`$. In order to find the asymptote we determine the integration constants $`C`$ and $`D`$ in Eq. (14) such that $`a_m(L_2)=h(L_2)`$ and $`a_m^{}(L_2)=h^{}(L_2)`$ where $`z=L_2`$ is the coordinate up to which the differential equation is integrated numerically; the system size $`L_2L_1`$ is chosen large enough so that upon further increase of the system size $`C`$ and $`D`$ remain unchanged within prescribed accuracy. The contact angle $`\theta _m`$ can be inferred from the value of $`C`$ (Eq. (17)). Finally the coordinate system is shifted laterally such that the intersection of $`a_m(z)`$ and $`h_0`$ define the position $`z=0`$ (which corresponds to $`z=z_1`$ for $`\kappa 0`$). As mentioned before, for macroscopic drops the rotationally symmetric configuration satisfies the stability criterion Eq. (23) for any $`R`$.
The dependence of the liquid-vapor interface profiles of macroscopic drops on the cylinder radius $`R`$ is shown in Fig. 8 using $`\omega _p(l)`$ and with the same set of interaction potential parameters as in the previous examples. In accordance with Eq. (17), the interface profiles for cylinders of finite thickness diverge exponentially in the limit $`z\mathrm{}`$. In the limit $`R\mathrm{}`$ the region where higher-order corrections to the linear behavior (Eq. (18)) are relevant is shifted towards $`z\mathrm{}`$ such that in the limiting case $`R=\mathrm{}`$ corresponding to the planar substrate the linear divergence of the reference profile is recovered. Figure 9(a) displays the apparent contact angles $`\theta _m`$ corresponding to the profiles shown in Fig. 8 as function of $`R`$. Upon increasing the cylinder radius $`R`$, $`\theta _m`$ approaches Young’s contact angle $`\theta _{\mathrm{}}`$ for the planar substrate as $`\theta _{\mathrm{}}\theta _mR^1`$.
In Fig. 9(b) we show the apparent contact angles of *small* droplets on a *planar* substrate, i.e., in the limit $`R\mathrm{}`$ but with $`V_{ex}<\mathrm{}`$. In this case the reference configuration is a spherical cap whose circular base has a radius $`z_1`$ (see Fig. 1(a)). For large droplets $`\mathrm{cos}\theta `$ reaches Young’s contact angle $`\mathrm{cos}\theta _{\mathrm{}}`$ according to the Neumann-Boruvka equation
$$\mathrm{cos}\theta \mathrm{cos}\theta _{\mathrm{}}=\frac{\tau _{\mathrm{}}}{\sigma z_1}$$
(24)
which allows one to determine experimentally the line tension $`\tau _{\mathrm{}}`$ of three-phase contact on a planar substrate by varying the droplet size. Figure 9(b) demonstrates that this linear relationship between $`\mathrm{cos}\theta `$ and $`z_1^1`$ is valid only for $`z_1/s500`$, i.e., for $`z_1500`$nm. From Fig. 9(b) one infers that $`\mathrm{cos}\theta `$ decreases more rapidly than predicted by Eq. (24). This behavior can be accounted for by an effective line tension $`\tau _{eff}(z_1)`$ which due to the circular bending of the three-phase contact line is *larger* than the value $`\tau _{\mathrm{}}`$ of the corresponding *straight* three-phase contact line. Similar results have been obtained by Dobbs .
$`\theta `$ and $`\theta _m`$ depend on the liquid-vapor surface tension $`\sigma `$ which in turn also exhibits a behavior $`\sigma (R)\sigma (\mathrm{})R^1`$. In accordance with the discussion in Subsec. II A, Fig. 9 does not yet take into account this indirect dependence of $`\theta `$ and $`\theta _m`$ on $`R`$ via $`\sigma (R)`$.
## IV Line tension
As long as the size of a droplet is *finite* and *fixed* it is impossible to extract from the total free energy unambiguously and in a strict thermodynamic sense a line tension associated with the three-phase contact lines at the ends of the droplet because there are arbitrarily many ways to form the total free energy as a sum of various terms. However, *well-defined* line tensions emerge as coefficients in the *size dependence* of the free energy of droplets upon approaching macroscopic drops. To this end we consider the limit of large drops, i.e., $`V_{ex}s^31`$ and $`z_1/s1`$ (see Fig. 1(a)). Within the interface displacement model the excess free energy in Eq. (5) can be rewritten as
$$F_{ex}=\sigma A_b4\pi R\omega _c(l_0)z_1\sigma A_c+$$
(25)
where
$$A_b=2\pi \underset{z_1}{\overset{z_1}{}}𝑑za_{ref}(z)\sqrt{1+(a_{ref}^{}(z))^2}$$
(26)
is the surface area of the “barrel” part of the reference surface $`a_{ref}(z)`$, $`z_1zz_1`$, $`A_c=4\pi h_0z_1`$ is the surface area of the cylinder with radius $`h_0`$ and length $`z_1`$ (see Fig. 1(a)), and
$``$ $`=`$ $`2\pi {\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}dz(\sigma (h(z)\sqrt{1+(h^{}(z))^2}a_{ref}(z)\sqrt{1+(a_{ref}^{}(z))^2})`$ (28)
$`+R(\omega (h(z)R)\omega (a_{ref}(z)R)))+2\pi R{\displaystyle }_{z_1}^{z_1}dz\omega (a_{ref}(z)R).`$
For $`V_{ex}\mathrm{}`$, $`F_{ex}`$ (Eq. (25)) is dominated by the term $`\sigma A_b`$ which scales proportional to the surface area of the drop and thus represents a two-dimensional contribution. The leading subdominant terms are $`4\pi R\omega _c(l_0)z_1`$ and $`4\pi h_0\sigma z_1`$, which scale with the linear dimension $`2z_1`$ of the drop representing one-dimensional contributions. Finally, the last term in Eq. (25) remains finite for $`V_{ex}`$, $`A_b`$, and $`z_1\mathrm{}`$:
$$(z_1)=𝒯+𝒪(z_1^1)$$
(29)
and thus represents a zero-dimensional contribution. In its turn $`𝒯`$ depends on the cylinder radius $`R`$ such that for large $`R`$ it scales proportional to $`R`$ which leads to the following definition of an excess free energy per unit length, henceforth called “line tension”, associated with the two contact lines formed at the ends of the droplet with total length $`4\pi R`$:
$$\tau =\frac{𝒯}{4\pi R}.$$
(30)
We note that $`𝒯`$ is only well-defined in the thermodynamic limit $`\underset{A_b,z_1\mathrm{}}{lim}[F_{ex}\sigma A_b+4\pi R\omega _c(l_0)z_1+\sigma A_c]`$ so that higher order terms, e.g., $`z_1^1`$, omitted on the right hand side of Eq. (29) drop out. In the following $`𝒯`$ is understood to have been obtained via this procedure. It can be expressed in terms of the solution of Eq. (9) for $`\kappa =0`$ and of $`a_{ref,m}(z)`$ (Eqs. (14) and (15)):
$`𝒯`$ $`=`$ $`4\pi {\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}dz(\sigma (h(z)\sqrt{1+(h^{}(z))^2}a_{ref,m}(z)\sqrt{1+(a_{ref,m}^{}(z))^2})`$ (32)
$`+R(\omega (h(z)R)\omega (a_{ref,m}(z)R)))+4\pi R{\displaystyle }_0^{\mathrm{}}dz\omega (a_{ref,m}(z)R).`$
Consequently, $`𝒯`$ is twice the characteristic excess free energy associated with the structure of a macroscopic drop near one of its ends without interference from the other end. On the other hand, $`𝒯`$, and thus $`\tau `$, are defined for any value of $`R`$. The ratio $`\tau `$ formed in Eq. (30) has the property that in the limit $`R\mathrm{}`$ it reduces to the line tension $`\tau _{\mathrm{}}`$ of the straight three-phase contact line on the corresponding planar substrate (see, e.g., Refs. ), which is an experimentally observable quantity (compare Fig. 9(b)).
At this stage one should note that the above considerations tacitly assume that another thermodynamic limit concerning the total system size, such as the volume of the surrounding vapor phase and the length $`L`$ of the solid cylinder, has already been carried out in advance: $`F(h_0)`$ is proportional to $`L`$ and has been subtracted before. Moreover, we have not considered the bulk free energy of the surrounding vapor phase and the bulk free energy of the liquid in the drop proportional to $`V_{ex}`$ because they do not enter the description of the droplet shape in terms of an interface displacement model. As a careful analysis of the line tension $`\tau _{\mathrm{}}`$ within density functional theory for a volatile liquid at gas-liquid coexistence shows, in comparison with this more complete theory the interface displacement model misses a contribution which is independent of the shape $`l(x)`$ and is determined by $`\theta _{\mathrm{}}`$ and $`l_0`$ (the first term in the sum in Eq. (2.19) in Ref. , denoted as $`\stackrel{~}{\tau }`$ in Eqs. (4.2), (4.3), and (4.5) in Ref. ). Since, however, this constant contribution turns out to be numerically much smaller than those contributions captured by the interface displacement model (see Fig. 15 in Ref. ), we have refrained from determining it for the present, much more complicated geometry, assuming that the size ratios of these types of contributions remain roughly the same for the planar and the cylindrical substrate.
Figure 10 shows the dependence of the line tension $`\tau `$ on the cylinder radius $`R`$ using the planar effective interface potential $`\omega _p(l)`$. $`\tau `$ is given by $`s\sigma `$ times a numerical factor of the order of 1. It decreases monotonously for decreasing $`R`$ and attains its maximum value $`\tau _{\mathrm{}}`$ for $`R\mathrm{}`$ as $`\tau _{\mathrm{}}\tau (R)R^1`$. We note that this decrease of the line tension upon decreasing the radius of curvature $`R`$ of the contact line is opposite to the increase of the effective line tension $`\tau _{eff}(z_1)`$ observed for a decreasing radius of curvature $`z_1`$ of the circular three-phase contact line on a planar substrate as can be inferred from Fig. 9(b). Thus line tensions of curved three-phase contact lines can be smaller or larger than the line tension of the corresponding straight contact lines.
Whereas $`\tau _{\mathrm{}}`$ is experimentally accessible by monitoring the apparent contact angle of sessile droplets on a planar substrate as function of the droplet size, $`\tau (R)`$ cannot be determined experimentally by direct observation. The basic reason for this difference is that the length $`2\pi z_1`$ of the three-phase contact line of the sessile drop on the planar substrate can vary as function of the droplet size so that the optimal shape of the droplet responds to the associated cost $`2\pi z_1\tau _{\mathrm{}}`$ of the free energy, whereas the excess free energy $`2\pi R\tau (R)`$ for the ends of the droplets is a constant contribution with respect to the droplet size on the cylinder due to the fixed value of $`R`$. This, however, holds only for the “barrel”-type shape of the drop, for which the length of the three-phase contact lines is fixed. For “clamshell”-type droplet shapes the length of the three-phase contact line does depend on the droplet size so that in this case the line tension will influence the droplet shape. Nonetheless there are systems for which $`\tau (R)`$ can be experimentally relevant. If the cylinder is not a hard solid rod but consists of a soft material like, e.g., vesicles or tobacco viruses which float vertically at the liquid-vapor interface of a solvent, the positive line tension $`\tau (R)`$ will strangle the object locally, depending on its restoring elastic forces. According to Fig. 10 this tweaking force $`d(\tau (R)R)/dR`$ is weaker for thin cylinders.
## V Summary
Based on an interface displacement model (Eq. (5)) we have analyzed the shape and the free energy of “barrel”-type droplets of fixed volume $`V_{ex}`$ covering a cylindrical substrate of radius $`R`$ (Fig. 1). For sufficiently small droplets their shape $`h(z)=R+l(z)`$ is not only governed by the surface tension $`\sigma `$ of the liquid-vapor interface but also by the effective interface potential $`\omega _c(l;R)`$ with a generic form as shown in Fig. 2. We have obtained the following main results:
1. Figures 3 and 4 show how the deviation of the actual droplet shape $`h(z)`$ from a suitably defined reference configuration $`a_{ref}(z)`$ depends on the droplet size. The reference configuration is uniquely defined by the requirement to touch the actual shape at the apex and to have a constant mean curvature which equals the actual one at the apex. $`a_{ref}(z)`$ allows one to introduce an apparent contact angle $`\theta `$, characterizing the actual shape, which can be expressed in terms of the experimentally accessible quantities cylinder radius $`R`$, radii of curvature at the apex, height $`l(0)`$ of the droplet, and thickness $`l_0`$ of the wetting layer outside the barrel (Eq. (4)).
2. The dependence of the effective interface potential $`\omega _c(l;R)`$ on the cylinder radius $`R`$ influences the shape of the droplet on the scale of the range $`s`$ of $`\omega _c(l;R)`$ (Fig. 5); this dependence has a rather marked effect on the apparent contact angle (Fig. 7).
3. The apparent contact angles increase for smaller droplets and for thicker cylinders (Fig. 6). The contact angles of macroscopicly large drops approach Young’s contact angle on a planar substrate proportional to $`1/R`$ (Fig. 9(a)).
4. In the limiting case of small droplets on a planar substrate the circular bend of the three-phase contact line leads to an effectively increased value of the corresponding line tension (Fig. 9(b)). This deviation from the Neumann-Boruvka equation becomes relevant if the droplet radius is less than roughly $`500`$nm. This observation is relevant for experimental determinations of line tensions via contact angle measurements.
5. For macroscopicly large drops their shape $`l(z)`$ increases linearly on a planar substrate but exponentially on a cylinder (Fig. 1(b)). Figure 8 illustrates the smooth crossover between these types of behavior for increasing cylinder radii.
6. For large cylinder radii the line tension associated with the ends of macroscopicly large drops approaches the line tension of three-phase contact on a planar substrate proportional to $`1/R`$ (Fig. 10).
###### Acknowledgements.
This work has been supported by the German Science Foundation within the Special Research Initiative *Wetting and Structure Formation at Interfaces*. |
warning/0002/cond-mat0002378.html | ar5iv | text | # A first-order transition with power-law singularity in models with absorbing states
## I Introduction
Recently, nonequilibrium phase transitions have been intensively studied in variety of models . In addition to some potential applications, the motivation to study these transitions comes from the belief that they can be categorized into universality classes similarly to equilibrium phase transitions. In this context, models which exhibit transitions between active and absorbing phases are of particular interest. There already exists substantial numerical evidence that phase transitions in such models indeed can be classified into some universality classes. For example, it is believed that models with unique absorbing states should belong to the so-called directed-percolation (DP) universality class . Moreover, models with double (symmetric) absorbing states or with some conservation law in their dynamics, belong to another universality class .
Similarly to equilibrium systems, nonequilibrium continuous phase transitions are not the only possibility - some models are known to undergo discontinuous transitions . Although such transitions are not classified into universality classes, they might be more relevant since, discontinuous transitions are, at present, the only type of transitions which can be observed experimentally. On the contrary, the experimental realization of continuous phase transitions still remains an open problem .
One reason for a relatively good understanding of equilibrium phase transitions is a wealth of exactly solvable models in this field . With this respect, the situation is much worse for nonequilibrium phase transitions. None of the models with absorbing states and with continuous or discontinuous transitions was solved exactly and all results concerning the critical exponents or the location of a transition point are only numerical.
In the present paper we study certain models with infinitely many absorbing states. At a certain value of a control parameter $`r=r_\mathrm{c}`$, these models undergo a transition between active and absorbing phases. But the interesting point is a novel type of this transition: it seems to combine some features of both discontinuous and continuous transitions. Namely, at $`r=r_\mathrm{c}`$ an order parameter jumps discontinuously to zero but in addition to that the order parameter has a power-law singularity upon approaching the transition point from the active phase. Moreover, some elementary arguments, supported by Monte Carlo simulations, prompted us to predict the exact location of the transition point in both models, namely, $`r_\mathrm{c}=0`$. Both models have a gauge-like symmetry, which might be responsible for the unusual behaviour of these models.
## II Square lattice
Our first model is a certain modification of a model introduced in a context of modelling biological evolution . It is defined on a two-dimensional ($`d=2`$) square lattice where for each bond between the nearest-neighbouring sites $`i`$ and $`j`$ we introduce bond variables $`w_{i,j}(0.5,0.5)`$. Introducing the parameter $`r`$, we call the site $`i`$ active when $`_jw_{i,j}<r`$, where $`j`$ runs over all nearest neighbours of $`i`$. Otherwise, the site is called nonactive. The model is driven by random sequential dynamics and when the active site $`i`$ is selected, we assign anew, with uniform probability, four bond variables $`w_{i,j}`$, where $`j`$ is one of the nearest neighbours of $`i`$. Nonactive sites are not updated, but updating a certain (active) site might change the status of its neighbours.
An important quantity characterizing this model is the steady-state density of active sites $`\rho `$. How does $`\rho `$ change with the control parameter $`r`$? Of course, for $`r(0.5)^4=0.0625`$ all sites are active ($`\rho =1`$) for any distribution of bond variables $`w_{i,j}`$. It is natural to expect that for $`r<0.0625`$ and not too small there will be a certain fraction of active sites ($`\rho >0`$) and this fraction will decrease when $`r`$ decreases. Since the dynamical rules imply that the model has absorbing states with all sites nonactive ($`\rho =0`$), one can expect that at a certain $`r`$ the model undergoes a transition between the active and absorbing phases. On general grounds one expects that this transition might be either continuous and presumably of (2+1)DP universality class or discontinuous.
The existence of a transition is confirmed in Fig.1, which shows the density $`\rho `$ as a function of $`r`$ obtained using Monte Carlo simulations. The simulations were performed for the linear system size $`L=300`$ and we checked that the presented results are, within small statistical error, size-independent. After relaxing the random initial configuration for $`t_{\mathrm{rel}}=10^4`$, we made measurements during runs of $`t=10^5`$ (the unit of time is defined as a single on average update/lattice site). From this figure one can also see that the transition point $`r_\mathrm{c}`$ is located very close to $`r=0`$ and in the following we are going to show that it is very likely that in this model $`r_\mathrm{c}=0`$ (exactly).
First, we show that for $`r<0`$ the model is in the absorbing phase. The argument for that is elementary and based on the following observation: for $`r<0`$ there exists a finite probability that after updating a given site will become nonactive forever. Indeed, when one of the anew selected bonds satisfies the condition
$$|w_{i,j}|<r/(0.5)^3,$$
(1)
then the sites $`i`$ and $`j`$ become permanently nonactive (i.e., no matter what are the other bonds attached to these sites, they will always remain nonactive). For $`r<0`$ there is a finite probability to satisfy Eq. 1 and the above mechanism leads to the rapid decrease of active sites and hence the system reaches an absorbing state. The above mechanism is not effective for $`r0`$ since there is no value which would ensure permanent nonactivity of a certain site.
To confirm that for $`r<0`$ the system is in the absorbing phase, we present in Fig. 2 the time evolution of $`\rho `$ for $`r=10^6`$ and $`10^7`$. Although these values are very close to $`r=0`$, one can clearly see that the system evolves toward the absorbing state. (For $`r`$ smaller than these values, the approach to the absorbing state would be even faster.) As we have already mentioned, for $`r0`$ the mechanism which generates permanently nonactive sites is not effective. Most likely, this has important consequences: as shown in Fig. 2, even for $`r=0`$ the system does not evolve toward the absorbing state but remains in the active phase.
These results indicate that at $`r=0`$ the model undergoes a first-order transition between active and absorbing phases, characterized by a discontinuity of the order parameter $`\rho `$. However, the most interesting feature of the model is the fact that upon approaching the first-order transition point $`r=0`$ the order parameter exhibits a power-law singularity. Such singularities usually signals a continuous transition. This singularity, which is already visible in the inset of Fig. 1, is also presented in the logarithmic plot in Fig. 3. The parameter $`\rho _0=0.359`$ (i.e., the density of active sites for $`r=0`$) in Fig. 3 was obtained from the least-square analysis of small-$`r`$ ($`r10^3`$) data shown in Fig. 1 using the formula
$$\rho (r)=\rho _0+Ar^\beta ,$$
(2)
where we assumed that the critical point is located at $`r=0`$ . From the slope of the data in Fig 3, we estimate $`\beta =0.58(1)`$, which might suggest that the exponent $`\beta `$ for that model is the same as in the (2+1)DP . However, a characteristic feature of models of the DP universality class is that at the transition point the model falls into an absorbing state. Our model at the transition point ($`r=0`$) is not in the absorbing phase (see Fig. 2), but it enters the absorbing phase as soon as $`r`$ becomes negative. In addition, scaling behaviour of our numerical data persists on a relatively small interval of $`r`$ and asymptotically a different behaviour might sets in. Further arguments against DP criticality of this model are presented in the next section.
Let us notice that the above model is characterized by very large gauge-like symmetry. Indeed, inverting ($`\pm `$) four bond variables around any elementary square does not change the activity of sites. The gauge symmetry was examined for many equilibrium lattice models . However, the up-to-now examined models with absorbing states do not possess this kind of symmetry. It would be interesting to check whether the unusual properties of this model are related with this symmetry. In the following we examine a one-dimensional model which possesses a similar symmetry.
## III Triangular ladder
Let us examine a model defined on a one-dimensional ($`d=1`$) ladder-like lattice, where each site also has four neighbours (see Fig. 4). When defined with the same dynamical rules as the model examined in the previous section, this $`d=1`$ model also has an analogous gauge symmetry (see Fig. 4).
We examined the properties of this model using the same Monte Carlo method. Results of our simulations for the steady-state density $`\rho `$ are shown in Fig. 1 and Fig. 3. As is usually the case of models with absorbing states, Monte Carlo simulations of the $`d=1`$ version are much more accurate. For example, close to and at the transition point $`r=0`$ we simulated the system of the size $`L=310^5`$ and the simulation time was typically $`t=10^6`$. As a result we were able to probe a much closer vicinity of the transition point.
Our results indicate that the behaviour of $`d=1`$ and $`d=2`$ versions of this model is very similar. Both models exhibit a qualitatively the same transition at $`r=0`$. In the $`d=1`$ case our estimations of the critical parameters are: $`\rho _0=0.314827`$ and $`\beta =0.66(3)`$. A relatively good scaling behaviour in this case is confirmed on over two decades (see Fig. 3). The obtained value of the exponent $`\beta `$ clearly excludes the DP value (in the case of one-dimensional DP $`\beta =0.276486`$ .
To get further insight into the nature of the transition point we examined the size dependence of the relaxation time $`\tau `$. We measured the time needed for the system starting from the random initial configuration to reach the steady-state. Typically, at the continuous transition $`\tau `$ diverges as $`L^z`$ where $`z`$ is a positive exponent. For the one- and two-dimensional DP $`z=1.58`$ and 1.76, respectively. At the discontinuous transition one expects that $`\tau `$ remains finite in the thermodynamic limit (i.e., $`z=0`$). For $`r=0`$ and $`d=1`$ results of our measurements, shown in Fig. 5 are, in our opinion, inconclusive. They may suggest a power-law divergence with a small exponent $`z(0.2`$), but positive curvature of our data might asymptotically lead to $`z=0`$. On the other hand even if $`z=0`$, it is not certain whether $`\tau `$ remains finite or diverges, but slower than a power of $`L`$. For $`r>0`$ (i.e., off-criticality) the numerical data are similar to the $`r=0`$, but on general grounds one expects that $`\tau `$ remains finite in the thermodynamic limit.
## IV Summary
In the present paper we studied two models which exhibit remarkably similar and unusual behaviour. These models have a transition point which, although mainly of discontinuous nature (jump of the order parameter and $`z=0`$) has a certain feature of continuous transitions (power-law singularity of the order parameter).
The main weakness of our paper is the lack of any theoretical argument which would explain the behaviour of these models. Both models possess certain gauge-like symmetry. The role of such symmetry in models with absorbing states was not yet explored and it is possible that the behaviour of these models might be related with this symmetry.
Are there any indications that such transitions might take place in real systems? In our opinion, one of the possible applications might be related with phase transitions in nuclear physics. Indeed, there are some indications that multifragmentation of heavy nuclei resembles a phase transition which has both first- and second-order features . Such systems have been already modeled using Ising-like models. However, such an approach implicitly assumes a thermalization of the system, which is not obvious in these multifragmentation processes. Models with absorbing states might provide an alternative description of such processes.
## ACKNOWLEDGMENTS
I thank prof. Des Johnston for his hospitality and the Department of Mathematics of the Heriot-Watt University (Edinburgh, Scotland) for allocating computer time. I also thank H. Hinrichsen for interesting discussion. |
warning/0002/hep-th0002070.html | ar5iv | text | # Initial Conditions and the Structure of the Singularity in Pre-Big-Bang Cosmology
## 1 Introduction
The low energy effective equations of string theory provide cosmological solutions which might be applicable just below the string scale in the very early universe. In the pre-big-bang (PBB) scenario, suggested naturally by the spirit and the symmetries of Superstring theory, the universe starts in a low curvature, low coupling regime and then enters a stage of dilaton driven kinetic inflation . To address one of the main problems of cosmology, namely the problem of the initial conditions, this interesting picture has been developed further in , where the authors suggest that the initial state of the universe could have consisted of a bath of gravitational and dilatonic waves, some of which would have collapsed leading to the birth of a baby inflationary universe. These PBB bubble universes would give rise, finally, after a yet-to-be clarified graceful exit mechanism, to the observed Friedman-Robertson-Walker (FRW) world.
The main purpose of this paper is to develop a modified realization of the PBB bubble picture of Buonanno, Damour and Veneziano , in which the spherically symmetric collapse leading to inflationary PBB solutions is substituted by the interaction of strictly plane waves. This modification affects only the initial state of the universe, while near the (spacelike) caustic singularity the model shows similar behaviour to that discussed in , leading to Kasner-like structure. Representing exact solutions to the classical string equations of motion to all orders in the inverse string tension , it looks rather natural to modify the PBB picture by incorporating plane waves into the postulate of “asymptotic past triviality” . Moreover, this picture is attractive not only due to the exactness of plane wave backgrounds for string propagation, but most importantly, because of mutual “fatal attraction” exercised by the plane waves which leads to an inevitable gravitational collapse independently of their strength, unlike in the spherical picture.
In the scenario we propose, one starts with a model universe in a low coupling, low curvature regime with plane gravitational and matter waves which eventually will gravitationally interact <sup>4</sup><sup>4</sup>4In the gravitational sector we limit ourselves to the discussion of constantly polarized waves (diagonal metrics) on the grounds that in string theory the collision of gravitational waves with variable polarization may be generically mapped via T-duality into another problem where the variable polarization of the incoming waves is transformed into an non-vanishing value of the $`B_{\mu \nu }`$ field .. Colliding plane wave space-times have been investigated in detail in general relativity (see and references therein). A generic feature of the interaction of two plane waves is the formation of a strong space-like curvature singularity in the future . In the context of the PBB scenario, this singularity can be re-interpreted as an ordinary cosmological singularity.
The approach to the singularity from the past occurs through a Kasner-like behaviour, and this, in turn, can be analytically related to the initial data. This will provide the framework for a future quantitative study of the important problem of fine tuning of initial conditions leading to inflationary behaviour as $`t0^{}`$. This problem is well defined in our picture, unlike in the spherically symmetric case where a similar analysis does not seem to be possible. In the case at hand, the problem is a direct generalisation of the problem of determining Kasner exponents in scattering of pure plane gravitational waves with constant polarization . As far as the technical part of this paper is concerned, some already well-known results from general relativity and mathematical cosmology will be used and re-interpreted in a new light.
The whole question of naturalness of initial conditions is far from being settled and is under current discussion in the literature. The authors of have found that the solution of the flatness problem requires the introduction of two huge dimensionless parameters. On the other hand in it was argued that this might not constitute such a fine tuning if the initial scale is taken to be the whole horizon and not just the Planck scale. From a different point of view the genericity of PBB cosmology has been addressed in , showing that plane waves might be considered a generic initial state of PBB cosmology. In this paper, however, we will not enter into the discussion of whether fine tuning is needed in order to solve the flatness problem, but will rather stick to the possibility of inflation, leaving the resolution of these other questions to future work.
The paper is organized as follows. In Sec. 2 we review briefly the initial value problem for the collision of dilatonic and gravitational waves. In particular we provide a closed expression relating the Kasner exponents which characterize the asymptotic geometry near the caustic singularity with the initial conditions for the metric functions and the dilaton. In Sec. 3 we apply these results to investigate the range of initial conditions leading to PBB inflation, and whether these conditions are naturally met in the collision of plane gravitational waves. In Sec. 4 we study two particular geometries in the interaction region, the Nappi-Witten solution , and a family of Kantowski-Sachs metrics studied in . Finally, in Sec. 5 we will use thermodynamical considerations to argue that plane waves are good candidates to represent the primordial PBB universe, summarize our conclusions and indicate possible future directions of research.
## 2 Colliding plane waves with aligned polarization
One of the basic assumptions of the PBB scenario is the so-called “asymptotic past triviality” (APT) hypothesis. According to it, the universe starts in the asymptotic past in a low curvature and low string coupling regime where the physics is accurately described in terms of tree level string theory. The effective dynamics of the long-range fields is thus governed by the leading terms of the effective low energy string action where both quantum and $`\alpha ^{}`$ corrections are ignored. In four dimensions this action is given in the string frame by
$$S=d^4x\sqrt{g}e^\varphi \left(R+g^{\alpha \beta }_\alpha \varphi _\beta \varphi \frac{1}{12}H^{\alpha \beta \gamma }H_{\alpha \beta \gamma }\right)$$
(1)
where the dilaton $`\varphi `$ and the antisymmetric tensor field strength $`H_{\alpha \beta \gamma }=_{[\alpha }B_{\beta \gamma ]}`$ are introduced. Furthermore, we will assume throughout that the extra six spatial dimensions are compactified in some internal appropriate manifold considered to be non-dynamical.
Applying the conformal transformation
$$g_{\alpha \beta }e^\varphi g_{\alpha \beta }.$$
(2)
the action can be written in the usual Einstein-Hilbert form (Einstein frame). In this frame the equations of motion are given by
$`R_{\mu \nu }{\displaystyle \frac{1}{2}}g_{\mu \nu }R`$ $`=`$ $`{}_{}{}^{(\varphi )}T_{\mu \nu }^{}+^{(H)}T_{\mu \nu }`$ (3)
$`_\mu \left[\mathrm{exp}(2\varphi )H^{\mu \nu \lambda }\right]`$ $`=`$ $`0`$ (4)
$`\mathrm{}\varphi +{\displaystyle \frac{1}{6}}e^{2\varphi }H_{\alpha \beta \gamma }H^{\alpha \beta \gamma }`$ $`=`$ $`0`$ (5)
where
$`{}_{}{}^{(\varphi )}T_{\mu \nu }^{}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(\varphi _{,\mu }\varphi _{,\nu }{\displaystyle \frac{1}{2}}g_{\mu \nu }g^{\alpha \beta }\varphi _{,\alpha }\varphi _{,\beta })`$ (6)
$`{}_{}{}^{(H)}T_{\mu \nu }^{}`$ $`=`$ $`{\displaystyle \frac{1}{12}}e^{2\varphi }\left(3H_{\mu \lambda \kappa }H_\nu ^{\lambda \kappa }{\displaystyle \frac{1}{2}}g_{\mu \nu }H_{\alpha \beta \gamma }H^{\alpha \beta \gamma }\right)`$ (7)
In four dimensions the antisymmetric tensor field strength can be written in terms of the pseudoscalar field, $`b`$, as follows
$`H^{\mu \nu \lambda }=e^{2\varphi }ϵ^{\rho \mu \nu \lambda }b_{,\rho },`$ (8)
and since the solutions including the axion $`b`$ can be obtained from pure dilaton solutions via a $`SL(2,IR)`$ transformation leaving invariant the Einstein frame metric, we will ignore this field and concentrate on gravi-dilaton system in what follows.
In the following two assumptions are required to hold for an asymptotically past trivial (APT) initial state:
* APT<sub>1</sub>: The string theory is weakly coupled, i.e. $`g=e^{\varphi /2}1.`$
* APT<sub>2</sub>: The curvature in string units is small.
APT<sub>1</sub> ensures that classical string theory is valid and string loop corrections to (1) can be ignored, whereas APT<sub>2</sub> means that $`\alpha ^{}`$ corrections to the same action are negligible. Our starting point will be to assume that in the asymptotic past the universe is in a trivial state characterized by gravitational waves propagating in a flat space-time. Eventually, these plane waves will collide giving rise to non-trivial geometries in the interaction region, and in particular, possibly, to the nucleation of PBB universes. Since plane waves are exact string vacua, APT<sub>2</sub> is automatically satisfied in the space-time regions before the interaction . However, it is not at all clear that the exact conformal invariance of the background also holds for all the possible solutions describing the interaction region. From physical considerations one would expect to have at least one of these solutions in the interaction region corresponding to an exact string background, and that this solution would smoothly match the incoming plane waves along the null boundaries of the interaction region (see below).
The space-times representing interactions of plane waves have two commuting Killing vectors $`\zeta _1`$, $`\zeta _2`$ and it is possible to chose a system of adapted coordinates $`(u,v,x,y)`$ in which $`\zeta _1=_x`$ and $`\zeta _2=_y`$, whereas the longitudinal coordinates $`(u,v)`$ are null. In the $`u`$-$`v`$ plane, the resulting space-time can be divided into four different regions (cf. Fig. 1) :
* Region I ($`u<0`$, $`v<0`$) is flat space-time, described by the usual Minkowski line element
$$ds_I^2=2dudv+dx^2+dy^2$$
and a constant dilaton.
* Region II ($`u0`$, $`v0`$) is incoming plane wave 1, described by
$`ds_{II}^2=2dudv+F_1^2(u)dx^2+G_1^2(u)dy^2,`$ (9)
and dilaton field $`\varphi _1(u)`$.
* Region III ($`u0`$, $`v0`$) is incoming plane wave 2, described by
$`ds_{III}^2=2dudv+F_2^2(v)dx^2+G_2^2(v)dy^2,`$ (10)
along with $`\varphi _2(v)`$.
* Region IV ($`u0`$, $`v0`$) is the interaction region, described by
$`ds_{IV}^2=2e^Fdudv+G(e^\psi dx^2+e^\psi dy^2),`$ (11)
where $`F(u,v)`$, $`G(u,v)`$ and $`\psi (u,v)`$, as well as the dilaton field $`\varphi (u,v)`$, are functions of both $`u`$ and $`v`$.
In this problem the initial data are most conveniently posed on the null surfaces $`N_1N_2`$, where $`N_1\{v=0,u0\}`$ and $`N_2\{u=0,v0\}`$, which is the boundary of the interaction region IV. In the interior of this region one of Einstein’s equations reads
$`G_{uv}=0`$ (12)
which is solved by
$$G=a(u)+b(v)=1(\alpha u)^n(\beta v)^m.$$
Here $`\alpha ^1`$ and $`\beta ^1`$ are arbitrary positive length scales fixing the focal lengths of the incoming waves that we will set to 1 in the following. On the other hand, the integers $`n`$ and $`m`$ are determined by the boundary conditions.
Introducing two new coordinates $`\xi `$ and $`z`$ defined by ($`\alpha =\beta =1`$)
$`\xi `$ $``$ $`a(u)+b(v)=1u^nv^m`$ (13)
$`z`$ $``$ $`a(u)b(v)=u^nv^m`$ (14)
the metric (11) in the interaction region takes the familiar Einstein-Rosen form
$`ds^2=e^f(d\xi ^2+dz^2)+\xi (e^\psi dx^2+e^\psi dy^2).`$ (15)
The equations of motion for the metric functions and the dilaton are given by
$`\ddot{\psi }`$ $`+`$ $`{\displaystyle \frac{1}{\xi }}\dot{\psi }\psi ^{\prime \prime }=0`$ (16)
$`\ddot{\varphi }`$ $`+`$ $`{\displaystyle \frac{1}{\xi }}\dot{\varphi }\varphi ^{\prime \prime }=0`$ (17)
$`\dot{f}`$ $`=`$ $`{\displaystyle \frac{1}{2\xi }}+{\displaystyle \frac{\xi }{2}}\left(\dot{\psi }^2+\psi ^2\right)+{\displaystyle \frac{\xi }{2}}\left(\dot{\varphi }^2+\varphi ^2\right)`$ (18)
$`f^{}`$ $`=`$ $`\xi \dot{\psi }\psi ^{}+\xi \dot{\varphi }\varphi ^{}`$ (19)
where a dot and a prime denote differentiation with respect to $`\xi `$ and $`z`$ respectively. The equations for $`\psi `$ and $`\varphi `$ can be solved in terms of Bessel and Neumann functions by
$`V=k\mathrm{ln}\xi +\{A_\omega \mathrm{cos}[\omega (z+z_0)]J_0(\omega \xi )\}+\{B_\omega \mathrm{cos}[\omega (z+z_0)]N_0(\omega \xi )\}.`$ (20)
Here $`V`$ stands for either $`\psi `$ or $`\varphi `$ and $`\{\mathrm{}\}`$ denotes arbitrary linear combinations of the terms in curly brackets including those of the form $`_0^{\mathrm{}}A_\omega \mathrm{cos}[\omega (z+z_0)]J_0(\omega \xi )`$, $`_0^{\mathrm{}}B_\omega \mathrm{cos}[\omega (z+z_0)]N_0(\omega \xi )`$.
In order to relate the asymptotic behaviour near the singularity at $`\xi =0`$ to the initial data given on the boundary of the interaction region at $`\{(u,0)\}\{(0,v)\}`$ it is useful to introduce yet another set of coordinates $`r`$ and $`s`$ defined by
$$r\xi z,s\xi +z.$$
In this case the equations for $`\psi `$ and $`\varphi `$ take the form
$`\psi _{,rs}+{\displaystyle \frac{1}{2(r+s)}}(\psi _{,r}+\psi _{,s})`$ $`=`$ $`0`$ (21)
$`\varphi _{,rs}+{\displaystyle \frac{1}{2(r+s)}}(\varphi _{,r}+\varphi _{,s})`$ $`=`$ $`0.`$ (22)
These two equations, together with the initial data on the null boundaries of the interaction region $`\{\psi (r,1),\psi (1,s)\}`$ and $`\{\varphi (r,1),\varphi (1,s)\}`$ pose a well defined initial value problem. Both $`\psi (r,s)`$ and $`\varphi (r,s)`$ are $`C^1`$ (and piecewise $`C^2`$) functions. This problem was first solved by Szekeres in the case of pure gravitational waves. Here the notation of Yurtsever is used,
$`\psi (r,s)`$ $`=`$ $`{\displaystyle _1^s}𝑑s^{}\left[\psi _{,s^{}}(1,s^{})+{\displaystyle \frac{\psi (1,s^{})}{2(1+s^{})}}\right]\left[{\displaystyle \frac{1+s^{}}{r+s}}\right]^{\frac{1}{2}}𝒫_{\frac{1}{2}}\left[1+2{\displaystyle \frac{(1r)(s^{}s)}{(1+s^{})(r+s)}}\right]`$ (23)
$`+{\displaystyle _1^r}𝑑r^{}\left[\psi _{,r^{}}(r^{},1)+{\displaystyle \frac{\psi (r^{},1)}{2(1+r^{})}}\right]\left[{\displaystyle \frac{1+r^{}}{r+s}}\right]^{\frac{1}{2}}𝒫_{\frac{1}{2}}\left[1+2{\displaystyle \frac{(1s)(r^{}r)}{(1+r^{})(r+s)}}\right]`$
$`\varphi (r,s)`$ $`=`$ $`{\displaystyle _1^s}𝑑s^{}\left[\varphi _{,s^{}}(1,s^{})+{\displaystyle \frac{\varphi (1,s^{})}{2(1+s^{})}}\right]\left[{\displaystyle \frac{1+s^{}}{r+s}}\right]^{\frac{1}{2}}𝒫_{\frac{1}{2}}\left[1+2{\displaystyle \frac{(1r)(s^{}s)}{(1+s^{})(r+s)}}\right]`$ (24)
$`+{\displaystyle _1^r}𝑑r^{}\left[\varphi _{,r^{}}(r^{},1)+{\displaystyle \frac{\varphi (r^{},1)}{2(1+r^{})}}\right]\left[{\displaystyle \frac{1+r^{}}{r+s}}\right]^{\frac{1}{2}}𝒫_{\frac{1}{2}}\left[1+2{\displaystyle \frac{(1s)(r^{}r)}{(1+r^{})(r+s)}}\right]`$
where $`𝒫_{\frac{1}{2}}(x)`$ is a Legendre function. It is important to stress here that, in order to study the behaviour near the singularity, we are only concerned with the “nonzero mode” of the dilaton. Therefore, $`\varphi (r,s)`$ in Eq. (24) is normalized in such a way that $`\varphi (1,1)=0`$. However, we can always add an arbitrary constant to the dilaton field and still have a solution to the wave equation (22). In particular, we can tune the string coupling constant to small values in regions II and III (and of course in I as well) without affecting the structure of the singularity in region IV.
Near the singularity at $`\xi =0`$ Kasner behaviour is expected, and it is known that space-times admitting two abelian space-like Killing vectors with parallel polarization have an asymptotically velocity dominated singularity so that curvature effects become negligible there.
The nice feature of the colliding plane wave space-times is that the initial value problem is well posed and can be solved exactly. This allows to relate the Kasner exponents which describe the behaviour of the metric near the singularity to the initial data given on the null boundaries of the interaction region. In order to find this relationship we expand the Bessel functions around $`\xi =0`$, and proceeding along the lines of Yurtsever’s work write the following decomposition
$`\psi (\xi ,z)`$ $`=`$ $`ϵ(z)\mathrm{ln}\xi +d(z)+H(\xi ,z),`$
$`\varphi (\xi ,z)`$ $`=`$ $`\phi (z)\mathrm{ln}\xi +\stackrel{~}{d}(z)+\stackrel{~}{H}(\xi ,z),`$
where $`ϵ(z)`$, $`\phi (z)`$, $`d(z)`$ and $`\stackrel{~}{d}(z)`$ are independent of $`\xi `$, and $`H(\xi ,z)`$, $`\stackrel{~}{H}(\xi ,z)`$ vanish in the limit $`\xi 0`$ \[or $`r+s0`$ in $`(r,s)`$ coordinates\]. Thus, the Kasner exponents in this limit are determined entirely by the coefficients of the logarithmic terms in the above expansions.
These coefficients can be computed directly from Eqs. (23) and (24). For $`ϵ(z)`$ one finds
$`ϵ(z)`$ $`=`$ $`{\displaystyle \frac{1}{\pi \sqrt{1+z}}}{\displaystyle _z^1}𝑑s\left[(1+s)^{\frac{1}{2}}\psi (1,s)\right]_{,s}\left({\displaystyle \frac{s+1}{sz}}\right)^{\frac{1}{2}}`$ (25)
$`+`$ $`{\displaystyle \frac{1}{\pi \sqrt{1z}}}{\displaystyle _z^1}𝑑r\left[(1+r)^{\frac{1}{2}}\psi (r,1)\right]_{,r}\left({\displaystyle \frac{r+1}{r+z}}\right)^{\frac{1}{2}},`$
and a similar expression holds for the dilaton
$`\phi (z)`$ $`=`$ $`{\displaystyle \frac{1}{\pi \sqrt{1+z}}}{\displaystyle _z^1}𝑑s\left[(1+s)^{\frac{1}{2}}\varphi (1,s)\right]_{,s}\left({\displaystyle \frac{s+1}{sz}}\right)^{\frac{1}{2}}`$ (26)
$`+`$ $`{\displaystyle \frac{1}{\pi \sqrt{1z}}}{\displaystyle _z^1}𝑑r\left[(1+r)^{\frac{1}{2}}\varphi (r,1)\right]_{,r}\left({\displaystyle \frac{r+1}{r+z}}\right)^{\frac{1}{2}}.`$
By introducing the leading logarithmic behaviour of both $`\psi (\xi ,z)`$ and $`\varphi (\xi ,z)`$ into the equations for $`f(\xi ,z)`$, Eqs. (18) and (19), we readily get the solution for the metric function $`f(\xi ,z)`$ near the singularity at $`\xi =0`$ to be
$`f(\xi ,z){\displaystyle \frac{1}{2}}[ϵ^2(z)+\phi ^2(z)1]\mathrm{ln}\xi .`$ (27)
Hence the asymptotic behaviour of the metric when $`\xi 0`$ is given by
$`ds^2=\xi ^{a(z)}(d\xi ^2+dz^2)+\xi ^{1+ϵ(z)}dx^2+\xi ^{1ϵ(z)}dy^2`$ (28)
where $`a(z)\frac{1}{2}[ϵ^2(z)+\phi ^2(z)1]`$ (cf. also ).
Thus, we have completely specified the asymptotic behaviour of the metric near the caustic singularity in terms of the initial data for $`\psi (\xi ,z)`$ and $`\varphi (\xi ,z)`$, as encoded by the functions $`ϵ(z)`$ and $`\phi (z)`$. In the following section we will use this result to study the initial conditions in the collision problem leading to PBB inflation.
## 3 Conditions for pre-big-bang inflation
Now that the relation between the asymptotic form of the metric (28) and the initial conditions on the boundary of the interaction region is given, we can address the problem of determining what kind of initial data lead to PBB inflationary solutions. Transforming the solution (28) to the string frame and switching, once in the string frame, from conformal to cosmic time we find the following Kasner exponents (generically, these will be functions of $`z`$)
$`p_1(z)`$ $``$ $`{\displaystyle \frac{1+ϵ(z)+\phi (z)}{b(z)+2}}`$
$`p_2(z)`$ $``$ $`{\displaystyle \frac{1ϵ(z)+\phi (z)}{b(z)+2}}`$
$`p_3(z)`$ $``$ $`{\displaystyle \frac{b(z)}{b(z)+2}}`$ (29)
where $`b(z)\frac{1}{2}[ϵ^2(z)+\phi ^2(z)+2\phi (z)1]`$ and the subscripts 1,2,3 correspond to the $`x,y,z`$ directions respectively. These exponents satisfy the usual conditions for dilaton-vacuum Kasner solutions in the string frame
$$\underset{i=1}{\overset{3}{}}p_i(z)=1+\frac{2\phi (z)}{b(z)+2},\underset{i=1}{\overset{3}{}}p_i(z)^2=1.$$
The conditions for PBB inflation ($`p_1,p_2,p_3<0`$) are then translated into the following conditions on the functions $`b(z)`$, $`ϵ(z)`$ and $`\phi (z)`$:
$`2<b(z)<0`$ (30)
$`1ϵ(z)+\phi (z)<0`$ (31)
$`1+ϵ(z)+\phi (z)<0.`$ (32)
Actually, the first inequality, when expressed in terms of $`ϵ(z)`$ and $`\phi (z)`$, reads
$`ϵ(z)^2+[\phi (z)+1]^2<2`$ (33)
Thus, the set of points in the $`ϵ(z)`$-$`\phi (z)`$ plane for which we get PBB inflationary solutions near the singularity is the quadrant of the circle defined by (33) and bounded by the lines (31) and (32) as shown in Fig. 2. This quadrant is inscribed on the square defined by $`|ϵ(z)|<\sqrt{2}`$ and $`|\phi (z)+1|<\sqrt{2}`$.
Since for the inflating models $`\phi (z)<0`$,
$$g_{\mathrm{eff}}(\xi )^{\frac{1}{2}\phi (z)}+\mathrm{},(\xi 0^{})$$
so the effective string coupling constant diverges at the singularity. As a consequence, near $`\xi =0`$ quantum corrections will be large and will lead, hopefully, to regularization of the singularity and a graceful exit into the post-big-bang (i.e. FRW) phase.
Now the question of genericity of PBB inflation can be posed. In order to answer this we impose APT on the initial data. APT<sub>1</sub> demands that the string coupling constant has to be small. Since our equations are invariant under the shift of the dilaton by a constant $`\varphi \varphi +\mathrm{constant}`$ we can always fix this constant in such a way that the string theory is weakly coupled in regions I, II and III and thus also at the null boundaries $`N_1=\{v=0,u>0\}`$, $`N_2=\{u=0,v>0\}`$.
On the other hand, APT<sub>2</sub> is implemented by requiring small curvatures. We can mathematically express this condition by demanding the components of the Weyl tensor to be small on the initial null hypersurfaces $`N_1`$ and $`N_2`$. On $`N_1`$ the only non-vanishing component is $`\mathrm{\Psi }_4`$ and it is given by
$`\mathrm{\Psi }_4|_{N_1}={\displaystyle \frac{1}{2}}\psi _{uu}+{\displaystyle \frac{1}{4}}\psi _u\left[2{\displaystyle \frac{n1}{u}}+{\displaystyle \frac{3nu^{n1}}{1u^n}}{\displaystyle \frac{1u^n}{nu^{n1}}}(\psi _u^2+\varphi _u^2)\right]`$ (34)
where we have used the notation $`\psi _1(u)\psi (u,v=0)`$ and $`\varphi _1(u)\varphi (u,v=0)`$. On the other hand, on $`N_2`$ only $`\mathrm{\Psi }_0`$ is non-zero, and we find
$`\mathrm{\Psi }_0|_{N_2}={\displaystyle \frac{1}{2}}\psi _{vv}+{\displaystyle \frac{1}{4}}\psi _v\left[2{\displaystyle \frac{m1}{v}}+{\displaystyle \frac{3mv^{m1}}{1v^m}}{\displaystyle \frac{1v^m}{mv^{m1}}}(\psi _v^2+\varphi _v^2)\right]`$ (35)
where now<sup>5</sup><sup>5</sup>5In writing $`\mathrm{\Psi }_4`$ and $`\mathrm{\Psi }_0`$ we have set the original focal lengths of the incoming waves $`\alpha ^1`$, $`\beta ^1`$ to 1. We can restore these length scales by writing $`\mathrm{\Psi }_4|_{N_1}\alpha ^2\mathrm{\Psi }_4|_{N_1}`$ and $`\mathrm{\Psi }_0|_{N_2}\beta ^2\mathrm{\Psi }_0|_{N_2}`$ in Eqs. (34) and (35) respectively. Since the string length $`\mathrm{}_{st}`$ is the natural scale of the problem, we take $`\alpha ^1=\beta ^1=\mathrm{}_{st}=1`$ and measure all curvatures in string units. $`\psi _2(v)\psi (u=0,v)`$ and $`\varphi _2(v)\varphi (u=0,v)`$.
Thus, in order to satisfy generically the conditions of having low curvature in string units, $`\mathrm{\Psi }_4|_{N_1},\mathrm{\Psi }_0|_{N_2}1`$, we have to demand that all the derivatives of the metric functions $`\psi _1(u)`$, $`\psi _2(v)`$ that appear in expressions (34) and (35), as well as the corresponding derivatives of the dilaton field on the boundaries, $`\varphi _1(u)`$, $`\varphi _2(v)`$ are much smaller than $`1`$. From the condition $`\psi (u=0,v=0)=0`$ and $`\varphi (u=0,v=0)=0`$ \[in the latter case by $`\varphi (u,v)`$ we denote just the “nonzero mode” of the dilaton\] and the smallness of the derivatives we conclude that the functions $`\psi _1(u)`$, $`\psi _2(v)`$ as well as $`\varphi _1(u)`$ and $`\varphi _2(v)`$ are approximately constant and close to zero. Consequently, the initial data compatible with APT will satisfy
$`\psi (1,s)\mu _1=\mathrm{constant}1,\psi (r,1)\mu _2=\mathrm{constant}1,`$ (36)
and similar relation for the “nonzero mode” of the dilaton field
$`\varphi (1,s)\nu _1,\varphi (r,1)\nu _2`$ (37)
where again $`\nu _1`$ and $`\nu _2`$ are constants much smaller than 1. If we now make use of the expressions (25) and (26) that give us the gravitational and matter source functions $`ϵ(z)`$ and $`\phi (z)`$ in terms of the initial data, we find that $`ϵ(z)`$ is given by
$`ϵ(z)\left({\displaystyle \frac{1z}{1+z}}\right)^{\frac{1}{2}}{\displaystyle \frac{\mu _1}{\pi }}+\left({\displaystyle \frac{1+z}{1z}}\right)^{\frac{1}{2}}{\displaystyle \frac{\mu _2}{\pi }}.`$ (38)
This expression implies that $`ϵ(z)1`$ for a large range of values of $`z(1,1)`$, as long as $`\mu _1,\mu _21`$. On the other hand for the dilaton we get a similar relation
$`\phi (z)\left({\displaystyle \frac{1z}{1+z}}\right)^{\frac{1}{2}}{\displaystyle \frac{\nu _1}{\pi }}+\left({\displaystyle \frac{1+z}{1z}}\right)^{\frac{1}{2}}{\displaystyle \frac{\nu _2}{\pi }}`$ (39)
and again, since $`\nu _1,\nu _21`$, $`\phi (z)1`$ for a large range of values of $`z`$.
Thus, we have found that, on general grounds, APT selects the values of $`ϵ(z)`$-$`\phi (z)`$ in a region around the origin $`ϵ(z)=\phi (z)=0`$. If for $`ϵ(z)`$ this is consistent with the coordinates of those points in parameter space corresponding to the models for which the nucleation of PBB bubbles happens (see Fig. 2), in the case of $`\phi (z)`$ the situation is not that good, since to achieve PBB inflation we need $`1\sqrt{2}<\phi (z)<1`$. In any case, it is important to notice that the bounds imposed by APT are not equally strong for $`\psi (u,v)`$ and $`\varphi (u,v)`$. While in order to fulfill APT<sub>2</sub> we need both the first and second derivatives of $`\psi (u,v)`$ to be much smaller than 1 on $`N_1`$ and $`N_2`$, for $`\varphi (u,v)`$ we need just to demand this same condition on the square of the first derivative. Thus, APT is compatible with the hierarchy between the constants $`\mu _1,\mu _2`$ and $`\nu _1,\nu _2`$.
Nevertheless, the important conclusion we have reached is that, as a result of the gravitational wave collision, PBB inflation happens for a dense set of initial data, i.e. the PBB inflation becomes achievable in our scenario. This means that, once we have an inflationary solution, inflation is stable under small perturbations of the initial conditions that lead to small variations of the Kasner exponents.
## 4 Particular solutions in the interaction region
In the previous Section we have discussed the general asymptotic behaviour of the solutions near the curvature singularity and have shown that one may completely specify the structure of the singularity in terms of initial data posed on the null boundaries of the interaction region. In what follows, we discuss two particular examples of metrics in this region.
It is known that any metric with two commuting spatial Killing directions describes the interaction region in a colliding wave problem provided the appropriate boundary conditions are met. Here we will concentrate our attention on two cases which we think are of physical relevance. First, we will study, in the light of our approach, the solution of Nappi and Witten which describe an inhomogeneous universe with closed spatial sections of $`S^3`$ topology . The most interesting feature of this solution is the fact that it is an exact string background. After that we will consider dilatonic generalizations of the Schwarzschild metric, in order to make contact with the spherical collapse picture of Buonanno, Damour and Veneziano and the PBB inflating Kantowski-Sachs universes .
### 4.1 The Nappi-Witten cosmological solution
The four-dimensional cosmological model studied by Nappi and Witten results as the target space theory of a $`SL(2,IR)\times SU(2)/SO(1,1)\times U(1)`$ gauged Wess-Zumino-Witten model. The solution contains, besides the metric, non-trivial values for the dilaton and antisymmetric tensor field. Actually, the solution containing the non-vanishing B-field can be obtained by an $`O(2,2;IR)`$ rotation of the metric (for a review see )
$`ds^2=dt^2+dw^2+\mathrm{tan}^2wdx^2+\mathrm{cot}^2tdy^2,`$ (40)
together with the dilaton field
$`\varphi =\varphi _0\mathrm{log}(\mathrm{sin}^2t\mathrm{cos}^2w).`$ (41)
The above line element may be thought of as a product of two two-dimensional black holes with Euclidean and Lorentzian signatures, both being exact string backgrounds , corresponding to a $`SL(2,IR)/SO(1,1)\times SU(2)/U(1)`$ coset model.
In the Einstein frame (40) is given by
$`ds^2=e^{f(t,w)}(dt^2+dw^2)+K(t,w)[e^{\psi (t,w)}dx^2+e^{\psi (t,w)}dy^2]`$ (42)
with
$`f(t,w)`$ $`=`$ $`\mathrm{log}(1\mathrm{cos}2t)+\mathrm{log}(1+\mathrm{cos}2w)`$
$`K(t,w)`$ $`=`$ $`\mathrm{sin}2t\mathrm{sin}2w`$
$`\psi (t,w)`$ $`=`$ $`\mathrm{log}\mathrm{tan}t+\mathrm{log}\mathrm{tan}w,`$ (43)
the dilaton field being given by (41). To relate the Nappi-Witten solution with the plane wave collision problem, it is convenient to switch to a different set of coordinates, namely
$`\xi =\mathrm{sin}2t\mathrm{sin}2w,z=\mathrm{cos}2t\mathrm{cos}2w`$ (44)
in which the solitonic nature<sup>6</sup><sup>6</sup>6The name solitons is owed to the inverse scattering technique used to obtain these solutions rather than to their physical properties. It was realised later that in the diagonal case the gravi-soliton solutions are related to the well known Lamb-Rosen pulses , see for a review. Their role in the colliding wave problem was discussed in . of the solutions can be made explicit. Bakas also studied these solutions by applying the inverse scattering transform technique on a Kasner seed metric in a search to relate the Geroch group to the symmetries of the string theory. The above coordinate transformation can be inverted to give
$`\mathrm{log}\mathrm{tan}t`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(\mathrm{arc}\mathrm{cosh}{\displaystyle \frac{1+z}{\xi }}\mathrm{arc}\mathrm{cosh}{\displaystyle \frac{1z}{\xi }}\right),`$
$`\mathrm{log}\mathrm{tan}w`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(\mathrm{arc}\mathrm{cosh}{\displaystyle \frac{1+z}{\xi }}+\mathrm{arc}\mathrm{cosh}{\displaystyle \frac{1z}{\xi }}\right)`$ (45)
so in the new coordinates the metric function $`\psi (\xi ,z)`$ is given by
$$\psi (\xi ,z)=\mathrm{arc}\mathrm{cosh}\left(\frac{1+z}{\xi }\right)$$
whereas the dilaton becomes
$`\varphi (\xi ,z)=\mathrm{arc}\mathrm{cosh}\left({\displaystyle \frac{1z}{\xi }}\right)\mathrm{log}\xi .`$ (46)
In it was pointed out that the presence of at least two solitonic terms provide sufficient conditions for the continuity on the two different null boundaries of the region IV, if one is to interpret the spacetime in terms of plane wave interaction. This is indeed the case, since there is one soliton (the $`\mathrm{arc}\mathrm{cosh}`$ term) associated with the dilaton solution and another one in the transverse metric function $`\psi (\xi ,z)`$, and the contribution to the boundary condition of each of those is equivalent as if there where two solitons in the gravitational sector (for a general discussion of the boundary conditions in the plane wave collision problem see Ref. , in the string theory context see )
We can now express both $`\psi (\xi ,z)`$ and $`\varphi (\xi ,z)`$ in $`(r,s)`$-coordinates ($`r=\xi z`$, $`s=\xi +z`$), so the initial data for our problem on the boundaries $`N_1`$ and $`N_2`$ are specified by the functions
$`\psi (1,s)=0,\psi (r,1)=\mathrm{log}\left[{\displaystyle \frac{3r+2\sqrt{2(1r)}}{1+r}}\right].`$ (47)
For the dilaton, on the other hand, we find
$$\varphi (1,s)=2\mathrm{log}\left(1\sqrt{\frac{1s}{2}}\right),\varphi (r,1)=\mathrm{log}\left(\frac{1+r}{2}\right).$$
By substituting these expressions into (25) and (26) we may directly study the outcome of the collision near the singularity. We find that
$$ϵ(z)=1,\phi (z)=2.$$
It can be easily seen that these values of $`ϵ(z)`$ and $`\phi (z)`$ lie just on the boundary of the region of points for which the model undergoes PBB inflation. If we compute the Kasner exponents using (29) we find that
$$p_1=1,p_2=p_3=0,$$
so there is only one inflating direction, while the other two are “frozen”. Incidentally, the metric near the singularity corresponds to the T-dual of Milne space-time.
The original Nappi-Witten metric is obtained from the above solution by the $`O(2,2;IR)`$ rotation in string frame (followed by a rescaling of the $`x`$ coordinate, $`xBx`$, see ). Taking into account that $`O(2,2;IR)SL(2,IR)_\tau \times SL(2,IR)_\rho `$ the required transformation can be written as
$`\tau ^{}=\tau ,\rho ^{}={\displaystyle \frac{1}{\rho +B}},B0,`$ (48)
where $`\tau `$ and $`\rho `$ are the usual Kähler and complex structure moduli constructed from the string frame metric (see, for example, ). As discussed in the Einstein frame metric function $`\psi (\xi ,z)`$ remains invariant under the $`O(2,2;IR)`$ rotation so we have
$$\psi (\xi ,z)_{\mathrm{NW}}=\mathrm{arc}\mathrm{cosh}\left(\frac{1+z}{\xi }\right),$$
while for the new dilaton we find
$$\varphi (\xi ,z)_{\mathrm{NW}}=\varphi (\xi ,z)\mathrm{log}\left(B^2+\xi ^2e^{2\varphi (\xi ,z)}\right)$$
with $`\varphi (\xi ,z)`$ given by (46). In addition, we have a non-vanishing value for the B-field
$$B_{xy}(\xi ,z)=\frac{B}{B^2+\xi ^2e^{2\varphi (\xi ,z)}}.$$
Since $`\psi (\xi ,z)`$ is left unchanged by the rotation, the initial conditions for $`\psi (r,s)_{\mathrm{NW}}`$ on the null boundaries $`N_1`$, $`N_2`$ are again given by (47). Therefore, the gravitational source function $`ϵ(z)`$ remains invariant. On the other hand, the dilaton does transform under (48), so the initial conditions for the transformed dilaton are<sup>7</sup><sup>7</sup>7Notice again that in giving the initial conditions for the dilaton we are restricting to the “nonzero mode” defined by $`\varphi (1,1)=0`$ in $`(r,s)`$-coordinates. This means in particular that in the case of the Nappi-Witten solution we should write the arbitrary additive constant $`\varphi _0`$ in the dilaton as $`\varphi _0+\mathrm{log}(1+B^2)`$ in order to recover this nonzero mode when $`\varphi _0=0`$.
$`\varphi (1,s)_{\mathrm{NW}}`$ $`=`$ $`\mathrm{log}\left({\displaystyle \frac{1+r}{2}}\right),`$ (49)
$`\varphi (r,1)_{\mathrm{NW}}`$ $`=`$ $`\mathrm{log}\left[{\displaystyle \frac{3s}{2}}+{\displaystyle \frac{1B^2}{1+B^2}}\sqrt{2(1s)}\right].`$ (50)
Using Eq. (26) we can check that the scalar source function $`\phi (z)_{\mathrm{NW}}`$ vanishes. Consequently, the model lies outside the inflationary region in the $`ϵ(z)`$-$`\phi (z)`$ plane. If we evaluate the related Kasner exponents using (29) we get
$$p_2=1,p_1=p_3=0,$$
so the metric asymptotically approaches the Milne regime as $`\xi 0^{}`$.
In fact, we may perform a somewhat more general analysis; if we start with a solution characterized by some values of $`ϵ(z)`$, $`\phi (z)`$ within the inflationary region, after a generic $`SL(2,IR)_\rho O(2,2;IR)`$ rotation the resulting metric near the singularity will be characterized by the new functions
$`\overline{ϵ}(z)=ϵ(z),\overline{\phi }(z)=\phi (z)2.`$ (51)
In particular, for every model leading to PBB inflation we have $`\sqrt{2}1<\phi (z)<1`$, so the transformed function $`\overline{\phi }(z)`$ will satisfy $`\overline{\phi }(z)>1`$ and thus the metric will not inflate at the singularity. Since transformations in the $`SL(2,IR)_\rho `$ factor of $`O(2,2;IR)`$ are the ones generating background values of the B-field, one might be tempted to conclude that PBB inflation is not robust under the introduction of this field. On the other hand, some of the models which were not inflating before the transformation was performed, may happen to inflate after. Incidentally, all points in Fig. 2 with $`\phi (z)1`$ are preserved by transformations
$$\left(\begin{array}{cc}A& B\\ C& D\end{array}\right)SL(2,IR)_\rho $$
with $`D0`$, whereas they transform as in (51) when $`D=0`$.
### 4.2 Dilatonic Schwarzschild-like metric
In the original picture of Ref. , the nucleation of PBB bubbles comes through gravitational instability in the asymptotically trivial Universe. In our proposal, on the other hand, this nucleation is not so much due to gravitational instability of the gravitational wave gas, but rather, the result of the mutual focusing of these waves due to their nonlinear interaction. Needless to say that the initial conditions are specified by quite different initial data in both pictures. What we propose here is the closest thing one might think to represent the decomposition of the initial data into plane waves in a non-linear theory. In what follows we will consider a set of initial data expressed as plane waves producing the same behaviour in the interaction region, as if it were a particular case of spherically symmetric gravitational collapse and will relate this to the solutions discussed in and . The indication that this is possible relies on the previous studies where the solution first obtained by Ferrari and Ibáñez , and representing part of the black hole, were investigated in detail.
To this end, we start again with the Gowdy metric (42) specifying
$`f(t,w)`$ $`=`$ $`{\displaystyle \frac{1}{2}}[(a+1)^2+b^21]\mathrm{log}\mathrm{sin}2ta\mathrm{log}(1+\mathrm{cos}2t)`$ (52)
$`K(t,w)`$ $`=`$ $`\mathrm{sin}2t\mathrm{sin}2w`$ (53)
$`\psi (t,w)`$ $`=`$ $`a\mathrm{log}\mathrm{tan}t+\mathrm{log}\mathrm{sin}2t\mathrm{sin}2w`$ (54)
and dilaton field
$$\varphi (t,w)=\varphi _0+b\mathrm{log}\mathrm{tan}t$$
where the two constants $`a`$ and $`b`$ satisfy the condition $`a^2+b^2=4`$. A common feature of this uniparametric family of solutions is that they are spatially homogeneous and of Kantowski-Sachs type with positive spatial curvature<sup>8</sup><sup>8</sup>8This family of solutions corresponds to the family of closed Kantowski-Sachs cosmologies studied in , as can be seen by writing them in the coordinate system $`\tau =\frac{1}{2}(1+\mathrm{cos}2t)`$, $`\mathrm{x}=y`$, $`\phi =\pi +2x`$ and $`\theta =\pi +2w`$.. In particular, for $`a=2`$ and $`b=0`$, we obtain the “inside-horizon region” of a Schwarzschild black hole with $`4M^2=1`$. From (52)-(54) we see that the metrics are singular at $`t=0,\frac{\pi }{2}`$. Whenever $`b0`$ these are true curvature singularities with the curvature invariants blowing up. On the other hand, when $`b=0`$ ($`a=\pm 2`$) the apparent singularity at $`t=0`$ is just a coordinate singularity, while the one at $`t=\frac{\pi }{2}`$ remains a curvature singularity.
We can now rewrite these solutions using $`(\xi ,z)`$-coordinates defined by Eq. (44). Doing so the metric function $`\psi (\xi ,z)`$ and the dilaton field $`\varphi (\xi ,z)`$ are
$`\psi (\xi ,z)`$ $`=`$ $`{\displaystyle \frac{a}{2}}\left(\mathrm{arc}\mathrm{cosh}{\displaystyle \frac{1+z}{\xi }}+\mathrm{arc}\mathrm{cosh}{\displaystyle \frac{1z}{\xi }}\right)+\mathrm{log}\xi `$
$`\varphi (\xi ,z)`$ $`=`$ $`{\displaystyle \frac{b}{2}}\left(\mathrm{arc}\mathrm{cosh}{\displaystyle \frac{1+z}{\xi }}+\mathrm{arc}\mathrm{cosh}{\displaystyle \frac{1z}{\xi }}\right)`$
Changing into $`(r,s)`$-coordinates we readily get the initial conditions for $`\psi (r,s)`$ on the null boundaries $`N_1`$, $`N_2`$
$`\psi (1,s)`$ $`=`$ $`{\displaystyle \frac{a}{2}}\mathrm{log}\left[{\displaystyle \frac{3s+2\sqrt{2(1s)}}{1+s}}\right]+\mathrm{log}\left({\displaystyle \frac{1+s}{2}}\right),`$
$`\psi (r,1)`$ $`=`$ $`{\displaystyle \frac{a}{2}}\mathrm{log}\left[{\displaystyle \frac{3r+2\sqrt{2(1r)}}{1+r}}\right]+\mathrm{log}\left({\displaystyle \frac{1+r}{2}}\right).`$
For the dilaton field we find
$$\varphi (1,s)=\frac{b}{2}\mathrm{log}\left[\frac{3s+2\sqrt{2(1s)}}{1+s}\right],\varphi (r,1)=\frac{b}{2}\mathrm{log}\left[\frac{3r+2\sqrt{2(1r)}}{1+r}\right].$$
From these expressions we can evaluate $`ϵ(z)`$ and $`\phi (z)`$ to get
$$ϵ(z)=1a,\phi (z)=b$$
and since $`a`$ and $`b`$ satisfy $`a^2+b^2=4`$, we find that the values of $`ϵ(z)`$ and $`\phi (z)`$ lie on the circumference defined by
$`[ϵ(z)1]^2+\phi (z)^2=4.`$ (55)
In Fig. 3 we have plotted this curve in the $`ϵ(z)`$-$`\phi (z)`$ plane. We find that it crosses the region of points for which there is PBB inflation as $`\xi 0`$. Actually, the points of the circumference (55) within the inflationary region correspond to the set of models studied in for which both scale factors inflate (see Fig. 2 of Ref. ).
We have focused our attention above to the family of deformations of the Schwarzschild black hole labeled by a single parameter and in which homogeneity is preserved, i.e. the metric is of Kantowski-Sachs type. One may construct more general deformations of the Schwarzschild metric by considering higher-dimensional moduli spaces. Moreover, the family of Kantowski-Sachs solutions studied here can be extended to a two-parametric class of solutions with “homogeneous” longitudinal part of the metric, defined by
$`f(t,w)`$ $`=`$ $`{\displaystyle \frac{1}{2}}[(a_1+a_3)^2+(b_1+b_3)^21]\mathrm{log}\mathrm{sin}2t(a_1a_3+b_1b_3)\mathrm{log}\mathrm{sin}(1+\mathrm{cos}2t)`$
$`K(t,w)`$ $`=`$ $`\mathrm{sin}2t\mathrm{sin}2w`$
$`\psi (t,w)`$ $`=`$ $`a_1\mathrm{log}\mathrm{tan}t+a_2\mathrm{log}\mathrm{tan}w+a_3\mathrm{log}(\mathrm{sin}2t\mathrm{sin}2w)`$
with the dilaton field
$$\varphi (t,w)=\varphi _0+b_1\mathrm{log}\mathrm{tan}t+b_2\mathrm{log}\mathrm{tan}w+b_3\mathrm{log}(\mathrm{sin}2t\mathrm{sin}2w)$$
and the set of constants $`\{a_1,a_2,a_3\}`$ and $`\{b_1,b_2,b_3\}`$ satisfying the four conditions
$`(a_1+a_2)^2+(b_1+b_2)^2`$ $`=`$ $`4`$
$`(a_1a_2)^2+(b_1b_2)^2`$ $`=`$ $`4`$
$`(a_2+a_3)^2+(b_2+b_3)^2`$ $`=`$ $`1`$
$`a_2a_3+b_2b_3`$ $`=`$ $`0.`$
By solving these equations and studying the behaviour of the solution close to the singularity, we find that the resulting two-parametric family of models covers the region of the $`ϵ(z)`$-$`\phi (z)`$ plane defined by
$$1ϵ(z)^2+\phi (z)^29$$
which indeed contains the set of points for which PBB inflation occurs. However, if we want to relate this with the gravitational collapse picture of Ref. the solutions must possess a rotational symmetry and the only models in the family with an $`SO(3)`$ isometry are those with $`a_2=b_2=b_3=0`$, $`a_3=1`$, which precisely fall into the family of Kantowski-Sachs metrics that we have studied above.
## 5 Conclusions and outlook
In this paper we have proposed a picture where the PBB inflation is realized starting from a trivial asymptotic state. The idea is to start with strictly plane waves moving in different directions which interact gravitationally at some stage to produce a space-time singularity. The structure of the solutions close to the singularity is of the Kasner type, and we were able to relate analytically the initial data on the wave fronts to the Kasner exponents in the string frame. Although maybe the collapsing regions studied here are not as generic as for example those studied in , our picture has two basic advantages:
* The initial background is exact from the point of view of string propagation.
* One can completely determine the structure of the singularity in terms of the initial data provided by the incoming waves.
These two elements make of this scenario, at least, a solid test bench and a sort of theoretical laboratory for the PBB ideas.
We have seen that there exists a dense set of initial data leading to inflationary behaviour in the PBB phase and, therefore, there is a good chance for an inflationary universe to emerge during this phase. Since the Kasner exponents actually carry a space dependence, i.e. they depend in general on one coordinate $`z`$, different regions with different Kasner exponents experience different types of inflation. Therefore, although we start with the collision of two plane waves we obtain near the singularity a rich structure with the formation of different multiple PBB inflationary bubbles.
Although we have extracted these conclusions from a general analysis of the collision of two gravi-dilatonic waves, we have also studied two concrete examples of initial conditions leading to different geometries in the interaction region. The first one leads to a Nappi-Witten solution in region IV and it has the obvious interest of providing an exact string background also in this region. The second example corresponds to a spherically symmetric interaction region that may describe gravitational collapse as the result of the collision.
Another motivation to look at plane wave space-times as a probable initial state for the PBB universe may come from thermodynamical considerations. One expects the universe to start in the lowest possible state of gravitational entropy. Let us suppose, arguing in the spirit of Penrose’s hypothesis , that we relate the gravitational entropy to some homogeneous function constructed from all possible curvature invariants and that we normalize this function to be vanishing when the invariants vanish. With such a function at hand we formulate a kind of “Generalised Third Law of Thermodynamics” and assign zero gravitational entropy to those spacetimes for which all curvature invariants vanish. Interestingly, the FRW models do not fall into the class of zero entropy models according to our definition unlike in the Penrose’s case. To justify this, we argue that the lowest entropy state, apart from being the simplest, must be an exact string background. The space-times with all vanishing curvature invariants (plane waves) certainly do so, while the FRW universe does not.
The additional support to consider the gravitational entropy content of the plane wave geometry to be vanishing comes from yet a different, though not unrelated argument. One would usually tend to relate the gravitational entropy with the phenomenon of quantum particle creation. It is commonly accepted that quantum particle creation indicates whether a system is endowed with nontrivial gravitational entropy. Plane waves, due to their symmetry and to the fact that all the curvature invariants vanish, do not polarize the vacuum, so quantum particles are not created in the vicinity of plane waves . This is consistent then with the hypothesis of assigning zero gravitation entropy to the plane wave. Moreover, it looks as if time is not a player in the plane wave regime. Due to the absence of the global Cauchy surface, one may consider such a pure plane wave geometry as “timeless”. Therefore, until two such waves interact, no notion of time as defined by entropy change is appreciated. What happens further is beyond the scope of this paper.
One of the most interesting issues, untouched here, is the problem of whether the gravitational wave collision problem can be globally defined in terms of an exact string background. Provided one starts with initial states that are exact string backgrounds, what are the conditions for the data to evolve into the interaction region without breaking conformal invariance? We know that one such solution exists in this region, namely the Nappi-Witten solution that we studied in Sec. 4.1. Therefore, since we start with plane gravitational waves (which are exact string backgrounds) and we can make the transition over the null boundaries as much differentiable as we like, the question remains of whether this implies conformal invariance in the interaction region. On purely physical grounds one would be tempted to say that exact string backgrounds in the regions II and III of Fig. 1 will smoothly (i.e. $`C^{\mathrm{}}`$) extend to an exact background in region IV, at least if the full string equations share the uniqueness properties of the Einstein equations.
An interesting issue to address would be to refine the picture provided here to account for more realistic situations in which the primordial gravitational waves are not plane but localized. It has been argued that for “almost” plane gravitational waves singularities also occur as the result of their collision . In the case of “graviton beams”, there is also mutual focusing and maybe production of singularities that could serve as seeds for PBB bubbles.
Finally, there has been some work on the collision of plane waves at Planckian energies leading to black hole nucleation through tunneling. This kind of scenarios might be a way to find a semi-classical approximation of the formation of the singularity that could be applied to the graceful exit problem in PBB cosmology. This, and the thermodynamical ideas we have outlined above, we hope to be able to discuss in the future.
## Acknowledgements
We are grateful to Jacob Bekenstein for enlightening correspondence. Our special thanks go to Gabriele Veneziano for his valuable comments on the manuscript and interesting discussions. A.F. acknowledges the support of University of the Basque Country Grant UPV 122.310-EB150/98 and Spanish Science Ministry Grant PB96-0250. K.E.K. is supported by the Swiss National Science Foundation. The work of M.A.V.-M. has been supported by FOM (Fundamenteel Onderzoek van de Materie) Foundation and by University of the Basque Country Grants UPV 063.310-EB187/98 and UPV 172.310-G02/99, and Spanish Science Ministry Grant AEN99-0315. K.E.K. and M.A.V.-M. thank the Department of Theoretical Physics of The University of the Basque Country for hospitality. |
warning/0002/cond-mat0002077.html | ar5iv | text | # Modeling of Tunneling Spectroscopy in HTSC
## I INTRODUCTION
Tunneling measurements on HTSC have revealed a rich variety of properties and characteristics \[1-4\]. They may be classified according to their low and high energy features. With the low energy features we may attribute: (i) variable subgap shape of conductance, ranging from sharp, cusplike, to a flat, BCS-like feature ; (ii) voltage and temperature dependence of quasiparticle conductivity ; (iii) subgap structure ; (iv) zero bias conductance peak (ZBCP) . The high energy features include: (i) asymmetry of conductance peaks ,(ii) van Hove singularity (VHS), (iii) conductance shape outside of the gap region (background (BG)) and its asymmetry ; (iv) dip feature ; (v) hump feature .These features are collected in schematic Fig 1 .
While the tunneling spectroscopy on conventional superconductors allows directly to find the energy gap of superconductor, the same measurements in HTSC are not as easily interpreted. Some times the same experiments on the same samples show different results : cusplike or flat subgap feature, symmetric or asymmetric conductance peaks. Usually the sharpest gap features are obtained when BG is weakly decreasing. A quantitative measure of it is the ratio of the conductance peak height (PH) to the background conductance: PHB=PH/BG. When the BG conductance is decreasing, the $`PHB>2`$, but when BG conductance is linearly increasing ($``$V), $`PHB<2`$ . Kouznetsov and Coffey and Kirtly and Scalapino suggested that the linearly increasing BG is arising from inelastic tunneling. As was mentioned in , the conductance is dominated by quasiparticle tunneling and that the effect of Andreev reflection is not significant. A theoretical model for tunneling spectroscopy employing tight-binding band structure, $`d_{x^2y^2}`$ gap symmetry, group velocity and tunneling directionality was studied by Z. Yusof, J. F. Zasadzinski, L. Coffey and N. Miyahawa . An angle resolved photoemission spectroscopy (ARPES) band structure specific to optimally-doped BSCCO (Bi-2212) was used to calculate the tunneling density of states for a direct comparison to the experimental tunneling conductance. This model produces an asymmetric, decreasing conductance background, asymmetric conductance peaks and variable subgap shape, ranging from sharp, cusplike to a flat, BCS-like feature. A standard technique in analyzing the tunneling conductance is to use a smeared BCS function
$$N(E)=N(0)\frac{Ei\mathrm{\Gamma }}{\sqrt{(Ei\mathrm{\Gamma })^2\mathrm{\Delta }^2}}$$
(1)
in which a scattering rate parameter (lifetime broadening factor) $`\mathrm{\Gamma }`$ is used to take into account any broadening of the gap region in the DOS. Fig. 2 shows the DOS calculated by formula (1) at $`\mathrm{\Delta }=46`$ meV and $`\mathrm{\Gamma }=9`$ meV (a), $`\mathrm{\Gamma }=3`$ meV (b) and $`\mathrm{\Gamma }=0`$ (c). The characteristic features of the DOS is the flat subgap structure at small $`\mathrm{\Gamma }`$. This method can not explain the asymmetry of the conductance peaks observed in the tunneling experiments.
In the case of d-wave symmetry we have
Fig.2 s-wave DOS at $`\mathrm{\Delta }=46`$ meV and different values of $`\mathrm{\Gamma }`$, calculated by formula (1).
Fig.3 d-wave DOS at $`\mathrm{\Delta }=46`$ meV and different values of $`\mathrm{\Gamma }`$, calculated by formula (2).
$$N(E)=N(0)Re_0^{2\pi }\frac{d\varphi }{2\pi }\frac{Ei\mathrm{\Gamma }}{\sqrt{(Ei\mathrm{\Gamma })^2\mathrm{\Delta }_0^2cos^2(2\varphi )}}$$
(2)
and DOS calculated by this formula are presented in Fig. 3. The characteristic features of the DOS is the cusplike subgap structure. As was mentioned in this standard technique requires that the comparison be made with normalized tunneling conductance date, and since HTSC tunneling conductance can exhibit a varied and complex background shape, this procedure may ”filter out” too much information from the conductance data. An alternative is to simply normalize the data by a constant.
In the tunneling data were first normalized by constructing a ”normal state” conductance obtained by fitting the high bias data to a third order polynomial. The normalized conductance date were compared to a weighted momentum averaged d-wave DOS
$$N(E)=f(\varphi )\frac{Ei\mathrm{\Gamma }}{\sqrt{(Ei\mathrm{\Gamma })^2\mathrm{\Delta }_0^2cos^2(2\varphi )}}𝑑\varphi $$
(3)
Here $`f(\varphi )`$ is an angular weighting function, which allows for a better fit with the experimental date in the gap region. A weighting function $`f(\varphi )=1+0.4cos(4\varphi )`$ was used which imposed a preferential angular selection of the DOS along the absolute maximum of the d-wave gap and tapers off towards the nodes of the gap. This is a rather weak directional function since the minimum of $`f(0)`$ along the nodes of the d-wave gap is still none-negligible .
A. J. Fedro and D. Koelling have done the modeling of the normal state and superconducting DOS of HTSC, using tight-binding band structure, including the next nearest neighbors
$$\xi _k=2t(cos(k_xa)+cos(k_ya))+4t^{}cos(k_xa)cos(k_ya)\mu $$
(4)
The calculation showed two singularities in DOS: a van Hove singularity in the center of energy band due to saddle point near ($`\pi `$,0) at $`t^{}=0`$ and another at the lower edge of the energy band due to extra flattening out at (0,0). As extended s-wave and d-wave superconducting DOS were considered in case of hole-doped situation ($`\mu <0`$) for different hole concentration. The Fermi surface for $`t^{}=0`$ and $`t^{}=0.45t`$ at the same concentration, corresponding $`\mu /2t=0.187`$ which was used in are presented in Fig. 4a. It is needed to move up from the Fermi surface (set as the zero energy) to reach the point ($`\pi `$,0) in case $`t^{}=0`$ and move down in case $`t^{}=0.45t`$. So, for $`t^{}=0`$ the Fermi energy lies to the left of the van Hove singularity and will move away from it with increased hole doping while for $`t^{}=0.45t`$ it lies to the right and will move towards to with increased hole doping (See Fig. 4b where the DOS for $`t^{}=0`$ and $`t^{}=0.45`$ are presented ). For calculation of the superconducting DOS Fedro and Koelling used formula
$$N(E)=\frac{1}{2}\underset{k}{}(1+\frac{\xi _k}{E_k})\delta (EE_k)+(1\frac{\xi _k}{E_k})\delta (E+E_k)$$
(5)
This formula is the limit of the expression for tunneling density of states (6) at $`\mathrm{\Gamma }=0`$ and $`|T_k|^2=1`$, where $`T_k`$ is tunneling matrix element.
Fig.4 Fermi surfaces (left) and DOS (right) for $`t^{}`$=0 (solid lines) and $`t^{}`$=0.45t (dashed lines) in formula (4) at $`\frac{\mu }{2t}`$=-0.187 which corresponds hole - doped situation.
Fig.5 DOS for $`t^{}`$=0 at different $`\mathrm{\Gamma }`$ for s-wave symmetry (a) and d-wave symmetry (b), calculated by formula (6).
The Fig.5a shows the result of calculation of the DOS for $`t^{}`$ = 0 at $`\mathrm{\Gamma }=`$ 0.07, 0.1 and 0.2 meV for s-wave symmetry which reflect the results of Fedro and Koelling. Fig.5b shows the same DOS for d-wave symmetry. In both cases the Fermi energy (set as the zero of energy) lies to the left of the van Hove singularity. There is the peaks’ asymmetry which is more pronounced at large $`\mathrm{\Gamma }`$.
## II Models and Methods
In this paper we use the method for calculation of the DOS presented in . The tunneling DOS of a superconductor is determined by the imaginary part of the retarded single particle Green’s function
$$N(E)=\frac{1}{\pi }Im\underset{k}{}|T_k|^2G^R(k,E)$$
(6)
For the superconducting state
$$G^R(k,E)=\frac{u_k^2}{EE_k+i\mathrm{\Gamma }}+\frac{v_k^2}{E+E_k+i\mathrm{\Gamma }}$$
(7)
where $`u_k^2`$ and $`v_k^2`$ are the usual coherence factors,
$`u_k^2={\displaystyle \frac{1}{2}}(1+{\displaystyle \frac{\xi _k}{E_k}})`$ (8)
$`v_k^2={\displaystyle \frac{1}{2}}(1{\displaystyle \frac{\xi _k}{E_k}})`$ (9)
and $`\mathrm{\Gamma }`$ is the quasiparticle lifetime broadening factor. The energy spectrum of quasiparticles in the superconducting state is determined by
$$E_k=\sqrt{|\mathrm{\Delta }(k)|^2+\xi _k^2}$$
(10)
with the effective band structure extracted from ARPES experiments
$`\xi _k=C_0+0.5C_1[cos(k_xa)+cos(k_ya)]`$ (11)
$`+C_2cos(k_xa)cos(k_ya)+0.5C_3[cos(2k_xa)+cos(2k_ya)]`$ (12)
$`+0.5C_4[cos(2k_xa)cos(k_ya)+cos(k_xa)cos(2k_ya)]`$ (13)
$`+C_5cos(2k_xa)cos(2k_ya)`$ (14)
Here $`\xi _k`$ is measured with respect to the Fermi energy ($`\xi _k`$=0), and the phenomenological parameters are (in units of eV) $`C_0=0.1305`$, $`C_1=0.5951`$, $`C_2=0.1636`$, $`C_3=0.0519`$, $`C_4=0.1117`$, $`C_5=0.0510`$.
Fig.6(left) 3D-plot of energy spectrum of normal state according to formula (9).
Fig.7(center) 3D-plot of coherence factor $`u_k^2`$ according to formula (8).
Fig.8(right) Fermi surface corresponding to the $`\xi _k=0`$ in formula (10). Dark straight line shows the line of directional tunneling, the dashed lines show the angular spread $`\mathrm{\Theta }_0`$.
Fig. 6 shows the three dimensional image of function (10) . There are saddle point in ($`\pi `$,0) and flattening out of the energy band at (0,0) which lead to the van Hove singularities in the DOS. The three dimensional graph of the coherence factor $`u_k^2`$ is shown in Fig. 7.
Since quasiparticles with momentum perpendicular to the barrier interface have the highest probability of tunneling, the tunneling matrix element $`|T_k|^2`$ reveals a need for factors of directionality $`D(k)`$ and group velocity $`v_g(k)`$ . The group velocity factor is defined by
$$v_g(k)=|\stackrel{}{}_k\xi _k.\widehat{n}|=|\frac{\xi _k}{k_x}cos(\theta )+\frac{\xi _k}{k_y}sin(\theta )|$$
(15)
where the unit vector $`n`$ defines the tunneling direction as shown in Fig. 8, which is perpendicular to the plane of the junction.
The directionality function $`D(k)`$ is defined by
$$D(k)=exp[\frac{k^2(𝐤.\widehat{𝐧})^\mathrm{𝟐}}{(𝐤.\widehat{𝐧})^2\mathrm{\Theta }_0^2}]$$
(16)
Here $`\mathrm{\Theta }_0`$ defines the angular spread of the quasiparticle momentum with none-negligible tunneling probability with respect to n. The tunneling matrix element $`|T_k|^2`$ is written as
$$|T_k|^2=v_g(k)D(k)$$
(17)
The three dimensional graphs of group velocity $`v_g(k)`$, directionality $`D(k)`$ and tunneling matrix element $`|T_k|^2`$ functions are shown in Fig 9.
## III Results and discussions
Different factors may lead to the changing of the energy gap $`\mathrm{\Delta }_0`$ in HTSC. In particular, strong effects are caused by nonmagnetic impurities . In superconductors with d-wave symmetry the nonmagnetic impurities destroy the superconductivity very efficiently. Possibility to destroy of Cooper pairs by impurities leads to their finite lifetime. If the state with the quasiparticle is not stationary state, it must attenuate with time due to transitions to other states. The corresponding wave function has the form $`e^{i\xi (p)t/\mathrm{}\mathrm{\Gamma }t/\mathrm{}}`$, where $`\mathrm{\Gamma }`$ is proportional to the probability of the transitions to the other states. It may be interpreted as the energy of the quasiparticle has the imaginary addition $`i\mathrm{\Gamma }`$. The relation between $`\mathrm{\Gamma }`$ and
Fig.9 3D plot of the group velocity function (a), the directionality function (b) and the tunneling matrix element according to formulas (11), (12) and (13), correspondingly.
Fig.10 The changing of DOS with energy gap $`\mathrm{\Delta }_0`$ for s- (a,b) and d-wave (c,d) symmetry at $`\mathrm{\Gamma }`$=3 meV (a,c) and $`\mathrm{\Gamma }`$=9 meV (b,d) without of effects of directionality and group velocity.
lifetime of quasiparticle $`\tau _s`$ is $`\mathrm{\Gamma }=\mathrm{}/\tau _s`$. Hence, the impurities lead to changing $`\mathrm{\Delta }_0`$ and we may do modeling of the influence of impurities on tunneling conductance by numerical calculation of DOS N(E) considering different values of $`\mathrm{\Delta }_0`$ in formula (6). Here we present results of calculation N(E) at $`\mathrm{\Delta }_0=\alpha \mathrm{\Delta }_{00}`$, where $`\alpha `$=0.2, 0.4, 0.6, 0.8, 1 and $`\mathrm{\Delta }_{00}`$= 46 meV.
The peculiarities of the quasiparticle energy spectrum (10) play an essential role in explanation of the conductance features. Here, based on the numerical calculation of DOS we consider that the underlying asymmetry of the conductance peaks is primarily due to the features of quasiparticles energy spectrum. The d-wave gap symmetry simply enhances the degree of the peaks asymmetry. The last one is also changed by changing the tunneling direction.
Fig. 10 shows the results of the numerical calculations of the DOS at $`\mathrm{\Gamma }_0=3`$ meV (a,c) and $`\mathrm{\Gamma }_0=9`$ meV (b,d) without effects of group velocity and directionality as for s-wave (a,b) and d-wave (c,d) gap symmetry, correspondingly for different values of energy gap $`\mathrm{\Delta }_0`$. We have decreased the energy gap $`\mathrm{\Delta }_0`$, starting from $`\mathrm{\Delta }_0=46`$meV. For clarity we present only three characteristic curves, which corresponds to $`\alpha \mathrm{\Delta }_0`$ with $`\alpha =`$1, 0.6 and 0.2. We exclude the effects of group velocity and directionality to demonstrate that they are not responsible for peaks asymmetry.
There is the asymmetry of the quasiparticle peak heights as for s- and d-wave symmetry. So, the origin of the peaks asymmetry is not due to d-wave symmetry of the energy gap of HTSC. There is more flat subgap behavior of DOS in the case of s-wave symmetry in comparing with the d-wave case. The increase of lifetime broadening factor $`\mathrm{\Gamma }`$ leads to the enhance of the peaks’ asymmetry. There are van Hove singularities in the DOS at small $`\mathrm{\Gamma }`$. The increase of $`\mathrm{\Gamma }`$ leads to the confluence of the quasiparticle and VHS peaks and this results to the enhance of the DOS peaks asymmetry due to saddle point in energy spectrum (10) at ($`\pi `$,0). Also note to the asymmetry of the background as for s- and d-wave gap symmetry.
Fig. 11 shows the $`\mathrm{\Delta }_0`$-dependence of DOS taking into account the effects of group velocity and directionality at $`\mathrm{\Gamma }=3`$ meV (a,c) and $`\mathrm{\Gamma }=9`$ meV (b,d) for s-wave (a,b) and d-wave (c,d) gap symmetry. As in we have taken $`\mathrm{\Theta }=0.25`$ and $`\mathrm{\Theta }_0=0.1`$.
There is also quasiparticle peaks’ asymmetry similar to s- and d-wave cases. But in d-wave case the asymmetry is more stronger than in s-wave. The effects of group velocity and directionality lead to disappear of the VHS in DOS. The increase of $`\mathrm{\Gamma }`$ enhances the quasiparticle peak asymmetry. The most strong effect of energy band structure on the DOS occurs along $`k_x`$-axis due to van Hove singularity at $`(\pi ,0)`$.
Fig.11 The changing of DOS with energy gap $`\mathrm{\Delta }_0`$ for s- wave (a,b) and d-wave (c,d) gap symmetry at $`\mathrm{\Gamma }`$=3 meV (a,c) and $`\mathrm{\Gamma }`$=9 meV (b,d) with the effects of directionality and group velocity.
Fig.12 Effects of directionality on the DOS as in s-wave (a,b) and d-wave gap symmetry at $`\mathrm{\Gamma }`$=3 meV (a,c) and $`\mathrm{\Gamma }`$=9 meV (b,d).
Fig. 12 has demonstrated this effect. We have presented the DOS at different $`\mathrm{\Theta }`$ at $`\mathrm{\Gamma }=3`$ meV (a,c) and $`\mathrm{\Gamma }=9`$ meV (b,d) as for s-wave (a,b) and d-wave (c,d) gap symmetry. In the case of s-symmetry the position of the quasiparticle peaks is constant excluding the direction along $`k_x`$( $`\mathrm{\Theta }=0`$). We pay attention to the strong peaks’ asymmetry in this case.
In the case of d-wave symmetry we have practically the same behavior around $`k_x`$ direction as for s-wave, but energy gap is changed due to the $`\mathrm{\Theta }`$-dependence of $`\mathrm{\Delta }_0`$ and correspondingly, the quasiparticle peaks are shifted to the zero energy.
Fig.13 shows the $`\mathrm{\Theta }_0`$-changing of DOS at $`\mathrm{\Gamma }=3`$ meV (a,c) and $`\mathrm{\Gamma }=9`$ meV (b,d) as for s-wave (a,b) and d-wave (c,d) gap symmetry. The increase of $`\mathrm{\Theta }_0`$ means the taking into play (inclusion) the states, close to ($`\pi ,0`$). It is reflected as an appearance of the van Hove singularity as in case of s-wave and d-wave gap symmetry at small $`\mathrm{\Gamma }`$. The VHS is more pronounced in case of d-wave in comparing with s-wave symmetry. The increase of $`\mathrm{\Gamma }`$ leads to confluence of the quasiparticle and VHS peaks.
We consider that the absence VHS peak on the experimental $`dI/dV`$-characteristics means the enough large lifetime broadening factor $`\mathrm{\Gamma }`$ in that HTSC material.
The origin of the peaks asymmetry in the tunneling DOS was studied in by considering the role of the tunneling matrix element $`|T_k|^2`$ in the clean limit case ($`\mathrm{\Gamma }=0`$), where for the calculation N(E) was used the formula (5).
We repeat the explanation of the paper because we believe that the following conclusion on the origin of the peaks asymmetry must be different. At $`E>0`$ (positive bias voltages) the first term of (5) contributes to the N(E) because of $`\delta (E_kE)`$. In this case, as can see from Fig. 7 and Fig. 9c and Fig. 14a $`|T_k|^2`$ selects only a relatively short region of states in k-space in which $`u_k^2>0`$. These are the states with $`\xi _k>0`$ (above the FS). For the majority of states integrated over ( see again Fig.7 and Fig.9c). At $`E<0`$ (negative bias voltages) the second term of (5) contributes to the DOS because of $`\delta (E_k+E)`$. In this case, as can see from Fig.7 and Fig. 9c and Fig. 14b, $`|T_k|^2`$ selects out a large region of k states where $`v_k^2>0`$, in fact, equal to one. These states are below the Fermi surface, where $`\xi _k<0`$. The overall effect then is to have a large negative bias conductance compared to the positive one. This is true as for s- and d-wave symmetry. Hence, the underlying asymmetry of the conductance peaks is primarily due to the band structure $`\xi _k`$ and d-wave symmetry simply enhances the degree of asymmetry of the peaks. So, the peaks’ asymmetry existing as for s-wave and d-wave symmetry is sensitive to band structure $`\xi _k`$.
Fig.13 Numerical calculation of the quasiparticle DOS with a s-wave (a,b) and d-wave (c,d) gap symmetry at $`\mathrm{\Gamma }`$=3 meV (a,c) and $`\mathrm{\Gamma }`$=9 meV (b,d) for different spread $`\mathrm{\Theta }_0`$.
Fig.14 ARPES energy spectrum along $`\mathrm{\Theta }=0.25`$. The values of coherent factors correspond to $`E>0`$ (a) and $`E<0`$ (b).
In summary, by changing of the energy gap $`\mathrm{\Delta }_0`$ in HTSC one may model the influence of nonmagnetic impurities on the DOS. We consider that the asymmetry of the quasiparticle peaks is due to the specific features of the energy spectrum of HTSC and that the d-wave gap symmetry only enhances the peaks’ asymmetry.
The absence of the VHS peak on the experimental $`dI/dV`$-charactristics means the large enough lifetime broadening factor $`\mathrm{\Gamma }`$ in HTSC.
## IV Acknowledgement
We thank Professor Y. Sobouti and Professor M.R. H. Khajehpour for useful discussions. |
warning/0002/hep-ph0002231.html | ar5iv | text | # 𝐽/Ψ suppression in an expanding equilibrium hadron gas Work partially supported by the Polish Committee for Scientific Research under contract KBN - 2 P03B 030 18
## I Introduction
Since the paper of Matsui and Satz there is a steady interest in the problem of $`J/\mathrm{\Psi }`$ suppression in a heavy-ion collision. The question is if this suppression can be treated as a signature for a quark-gluon plasma or if it can be explained by $`J/\mathrm{\Psi }`$ absorption in a hadron gas which appears in the central rapidity region (CRR) of the collision .
In the following paper we shall continue our previous investigations into the problem of $`J/\mathrm{\Psi }`$ suppression observed in a heavy-ion collision (for experimental data see e.g. and references therein). Now, we shall focus on the dependence of the suppression on the initial energy density reached in the CRR.
In our model, $`J/\mathrm{\Psi }`$ suppression is the result of a $`c\overline{c}`$ state absorption in a dense hadronic matter through interactions of the type
$$c\overline{c}+hD+\overline{D}+X,$$
(1)
where $`h`$ denotes a hadron, $`D`$ is a charm meson and $`X`$ means a particle which is necessary to conserve the charge, baryon number or strangeness. The hadronic matter is in a state of an ideal gas of massive hadrons in thermal and chemical equilibrium and consists of all species up to $`\mathrm{\Omega }^{}`$ baryon. Time evolution is given here by conservation laws combined with assumptions about the space-time structure of the system. A corresponding equation of state of the ideal gas makes then possible to express gas parameters such as temperature and chemical potentials as functions of time.
An ideal gas of real hadrons has a very interesting feature: it cools much slower than a pion gas when expands longitudinally. We have checked numerically that for the initial energy densities $`ϵ_0`$ corresponding to initial temperatures $`T_0`$ of the order of 200 MeV and for the freeze-out $`T_{f.o.}`$ not lower than about 100 MeV, the time dependence of the temperature of the expanding gas still keeps the well-known form $`T(t)=T_0t^a`$ (we put $`t_0=1fm`$). Only the exponent $`a`$ changes from $`\frac{1}{3}`$ for massless pions to the values $`\frac{1}{5.6}`$-$`\frac{1}{5.3}`$ for massive realistic hadrons. As a result, the time of the freeze-out $`t_{f.o}`$ is much greater for the hadron gas than for the pion one. For instance, when we take $`T_0=200MeV`$ and $`T_{f.o.}=140MeV`$ we obtain $`t_{f.o.}=7.37fm`$ ($`a=\frac{1}{5.6}`$) for the hadron gas and $`t_{f.o.}=2.9fm`$ ($`a=\frac{1}{3}`$) for the pion gas. The lower $`T_{f.o.}`$, the stronger difference. For $`T_0=200MeV`$ and $`T_{f.o.}=100MeV`$ we have $`t_{f.o.}=48.5fm`$ ($`a=\frac{1}{5.6}`$) for the hadron gas and $`t_{f.o.}=8fm`$ ($`a=\frac{1}{3}`$) for the pion gas. This has a direct consequence for $`J/\mathrm{\Psi }`$ suppression: the longer the system lasts, the deeper suppression causes (see Fig. 1).
We are going to calculate a survival factor for $`J/\mathrm{\Psi }`$ when new, more realistic conditions are taken into account. We consider a hadronic gas which is produced in the CRR region. This gas expands both longitudinally and transversely. The longitudinal expansion is a traditional adiabatic hydrodynamical evolution , the transverse expansion is considered as the rarefaction wave. An initial energy density depends now on the impact parameter $`b`$ and on the geometry of the collision.
## II The expanding hadron gas
For an ideal hadron gas in thermal and chemical equilibrium, which consists of $`l`$ species of particles, energy density $`ϵ`$, baryon number density $`n_B`$, strangeness density $`n_S`$ and entropy density $`s`$ read ($`\mathrm{}=c=1`$ always)
$$ϵ=\frac{1}{2\pi ^2}\underset{i=1}{\overset{l}{}}(2s_i+1)_0^{\mathrm{}}\frac{dpp^2E_i}{\mathrm{exp}\left\{\frac{E_i\mu _i}{T}\right\}+g_i},$$
(3)
$$n_B=\frac{1}{2\pi ^2}\underset{i=1}{\overset{l}{}}(2s_i+1)_0^{\mathrm{}}\frac{dpp^2B_i}{\mathrm{exp}\left\{\frac{E_i\mu _i}{T}\right\}+g_i},$$
(4)
$$n_S=\frac{1}{2\pi ^2}\underset{i=1}{\overset{l}{}}(2s_i+1)_0^{\mathrm{}}\frac{dpp^2S_i}{\mathrm{exp}\left\{\frac{E_i\mu _i}{T}\right\}+g_i},$$
(5)
$$s=\frac{1}{6\pi ^2T^2}\underset{i=1}{\overset{l}{}}(2s_i+1)_0^{\mathrm{}}\frac{dpp^4}{E_i}\frac{(E_i\mu _i)\mathrm{exp}\left\{\frac{E_i\mu _i}{T}\right\}}{\left(\mathrm{exp}\left\{\frac{E_i\mu _i}{T}\right\}+g_i\right)^2},$$
(6)
where $`E_i=(m_i^2+p^2)^{1/2}`$ and $`m_i`$, $`B_i`$, $`S_i`$, $`\mu _i`$, $`s_i`$ and $`g_i`$ are the mass, baryon number, strangeness, chemical potential, spin and a statistical factor of specie $`i`$ respectively (we treat an antiparticle as a different specie). And $`\mu _i=B_i\mu _B+S_i\mu _S`$, where $`\mu _B`$ and $`\mu _S`$ are overall baryon number and strangeness chemical potentials respectively.
We shall work here within the usual timetable of the events in the CRR of a given ion collision (for more details see e.g. ). We fix $`t=0`$ at the moment of the maximal overlap of the nuclei. After half of the time the nuclei need to cross each other, matter appears in the CRR. We assume that soon thereafter matter thermalizes and this moment, $`t_0`$, is estimated at about 1 fm . Then matter starts to expand and cool and after reaching the freeze-out temperature it is no longer a thermodynamical system. We denote this moment as $`t_{f.o.}`$. As we have already mentioned in the introduction, this matter is the hadron gas, which consists of all hadrons up to $`\mathrm{\Omega }^{}`$ baryon. The expansion proceeds according to the relativistic hydrodynamics equations and for the longitudinal component we have the following solution (for details see e.g. )
$$s(t)=\frac{s_0t_0}{t},n_B(t)=\frac{n_B^0t_0}{t},$$
(7)
where $`s_0`$ and $`n_B^0`$ are initial densities of the entropy and the baryon number respectively. The superimposed transverse expansion has the form of the rarefaction wave moving radially inward with a sound velocity $`c_s`$ .
To obtain the time dependence of temperature and baryon number and strangeness chemical potentials one has to solve numerically equations (4 \- 6) with $`s`$, $`n_B`$ and $`n_S`$ given as time dependent quantities. For $`s(t)`$, $`n_B(t)`$ we have expressions (7) and $`n_S=0`$ since we put the overall strangeness equal to zero during all the evolution (for more details see ).
The sound velocity squared is given by $`c_s^2=\frac{P}{ϵ}`$ and can be evaluated numerically .
## III J/$`\mathrm{\Psi }`$ absorption in hadronic matter
In a high energy heavy-ion collision, charmonium states are produced mainly through gluon fusion and it takes place during the overlap of colliding nuclei. For the purpose of our model, we shall assume that all $`c\overline{c}`$ pairs are created at the moment $`t=0`$. Before the fusion, gluons can suffer multiple elastic scattering on nucleons and gain some additional transverse momentum in this way . This manifests for instance in the observed broadening of the $`p_T`$ distribution of $`J/\mathrm{\Psi }`$ . Following , we express this effect by the transverse momentum distribution of the charmonium states of the form
$$g(p_T,ϵ_0)=\frac{2p_T}{p_T^2_{J/\mathrm{\Psi }}^{AB}(ϵ_0)}\mathrm{exp}\left\{\frac{p_T^2}{p_T^2_{J/\mathrm{\Psi }}^{AB}(ϵ_0)}\right\},$$
(8)
where $`p_T^2_{J/\mathrm{\Psi }}^{AB}(ϵ_0)`$ is the mean squared transverse momentum of $`J/\mathrm{\Psi }`$ gained in an A-B collision with the initial energy density $`ϵ_0`$. The momentum can be expressed as (for details see )
$$p_T^2_{J/\mathrm{\Psi }}^{AB}(ϵ)=p_T^2_{J/\mathrm{\Psi }}^{pp}+Kϵ,$$
(9)
with $`K=0.27fm^3GeV`$ and $`p_T^2_{J/\mathrm{\Psi }}^{pp}=1.24GeV^2`$ taken from a fit to the $`J/\mathrm{\Psi }`$ data of NA38 Collaboration . The expression in (8) is normalized to unity and is treated as the initial momentum distribution of charmonium states here.
For the simplicity of our model, we shall assume that all charmonium states are completely formed and can be absorbed by the constituents of a surrounding medium from the moment of creation. It means that we neglect a whole complex process of $`J/\mathrm{\Psi }`$ formation as presented in . The main feature of the above-mentioned process is that, soon after the moment of production, the $`c\overline{c}`$ pair binds a soft gluon and creates a pre-resonance $`c\overline{c}g`$ state, from which, after a time of the order of 0.3 fm, a physical charmonium state is formed. This means that the possible nuclear absorption of charmonium is, in fact, the absorption of the $`c\overline{c}g`$ state. But the latest has the cross-section $`\sigma _{abs}=7.3mb`$, which is much higher than $`J/\mathrm{\Psi }Nucleon`$ absorption cross-section $`\sigma _{\psi N}35mb`$ obtained from p-A data . This justifies our assumption: taking into account $`c\overline{c}g`$ absorption instead of charmonium disintegration in the nuclear matter would only strengthen $`J/\mathrm{\Psi }`$ suppression.
According to the above assumption, charmonium states can be absorbed first in the nuclear matter and soon later, when the matter appears in the CRR, in the hadron gas. Since these two processes are separated in time, $`J/\mathrm{\Psi }`$ survival factor for a heavy-ion collision with the initial energy density $`ϵ_0`$, may be written in the form
$$𝒩(ϵ_0)=𝒩_{n.m.}(ϵ_0)𝒩_{h.g.}(ϵ_0),$$
(10)
where $`𝒩_{n.m.}(ϵ_0)`$ and $`𝒩_{h.g.}(ϵ_0)`$ are $`J/\mathrm{\Psi }`$ survival factors in the nuclear matter and the hadron gas, respectively. For $`𝒩_{n.m.}(ϵ_0)`$ we have the usual approximation
$${}_{n.m.}{}^{}(ϵ_0)\mathrm{exp}\{\sigma _{\psi N}\rho _0L\},$$
(11)
where $`\rho _0`$ is the nuclear matter density and $`L`$ the mean path length of the $`J/\mathrm{\Psi }`$ through the colliding nuclei. For the last quantity, we use the expression given in :
$$L(b)=\frac{1}{2\rho _0T_{AB}}d^2\stackrel{}{s}T_A(\stackrel{}{s})T_B(\stackrel{}{s}\stackrel{}{b})\left[T_A(\stackrel{}{s})+T_B(\stackrel{}{s}\stackrel{}{b})\right],$$
(12)
where $`T_{AB}(b)=d^2\stackrel{}{s}T_A(\stackrel{}{s})T_B(\stackrel{}{s}\stackrel{}{b})`$, $`T_A(\stackrel{}{s})=𝑑z\rho _A(\stackrel{}{s},z)`$ is the nuclear density profile function, $`\rho _A(\stackrel{}{s},z)`$ the nuclear matter density distribution and $`b`$ the impact parameter. How to obtain $`ϵ_0`$ as a function of $`b`$ will be presented further.
To estimate $`𝒩_{h.g.}(ϵ_0)`$ we follow the description presented in . We shall focus on the plane $`z=0`$ ($`z`$ is a collision axis) and put $`J/\mathrm{\Psi }`$ longitudinal momentum equal to zero. Now the $`p_T`$-dependent $`J/\mathrm{\Psi }`$ survival factor $`𝒩_{h.g.}(p_T)`$ is given by (for details see )
$$𝒩_{h.g.}(p_T)=d^2\stackrel{}{s}f_0(s,p_T)\mathrm{exp}\left\{_{t_0}^{t_f}𝑑t\underset{i=1}{\overset{l}{}}\frac{d^3\stackrel{}{q}}{(2\pi )^3}f_i(\stackrel{}{q},t)\sigma _iv_{rel,i}\frac{p_\nu q_i^\nu }{EE_i^{}}\right\},$$
(13)
where the sum in the power is over all taken species of scatters (hadrons), $`p^\nu =(E,\stackrel{}{p}_T)`$ and $`q_i^\nu =(E_i^{},\stackrel{}{q})`$ are four momenta of $`J/\mathrm{\Psi }`$ and hadron specie $`i`$ respectively, $`\stackrel{}{v}=\stackrel{}{p}_T/E`$ is the velocity of the former, $`\sigma _i`$ states for the absorption cross-section of $`J/\mathrm{\Psi }h_i`$ scattering and $`v_{rel,i}`$ is the relative velocity of $`h_i`$ hadron with respect to $`J/\mathrm{\Psi }`$. When $`M`$ denotes $`J/\mathrm{\Psi }`$ mass, $`M=3097`$ MeV, $`v_{rel,i}`$ reads
$$v_{rel,i}=\left(1\frac{m_i^2M^2}{(p_\nu q_i^\nu )^2}\right)^{\frac{1}{2}}.$$
(14)
The upper limit of the time integral in (13) , $`t_f`$, is equal to $`t_{f.o.}`$ or to $`t_{esc}`$ – the moment of leaving by a given $`J/\mathrm{\Psi }`$ of the hadron medium, if the final-size effects are considered and $`t_{esc}<t_{f.o.}`$. For $`\sigma _i`$ we have assumed that it equals zero for $`(p^\nu +q_i^\nu )^2<(2m_D+m_X)^2`$ and is constant elsewhere ($`m_D`$ is a charm meson mass, $`m_D=1867`$ MeV). For hadron specie $`i`$ we have usual Bose-Einstein or Fermi-Dirac distribution (we neglect any possible spatial dependence here)
$$f_i(\stackrel{}{q},t)=f_i(q,t)=\frac{2s_i+1}{\mathrm{exp}\left\{\frac{E_i^{}\mu _i(t)}{T(t)}\right\}+g_i}.$$
(15)
In the following, we shall consider only $`J/\mathrm{\Psi }`$ initial distribution $`f_0(s,p_T)`$ that factorizes into $`f_0(s)g(p_T)`$ and the momentum distribution $`g(p_T)`$ will be given by (8) . We assume at the first step that the transverse size of the hadron medium is much greater than $`t_{f.o.}`$ and also much greater than the size of the area where $`f_0(s)`$ is non-zero. Additionally we assume that $`f_0(s)`$ is uniform and normalized to unity. Note that the first assumption overestimates the suppression but the second, in the presence of the first, has no any calculable effect here. As a result, $`𝒩_{h.g.}(p_T)`$ simplifies to
$$𝒩_{h.g.}(p_T)=g(p_T,ϵ_0)\mathrm{exp}\left\{_{t_0}^{t_{f.o.}}𝑑t\underset{i=1}{\overset{l}{}}\frac{d^3\stackrel{}{q}}{(2\pi )^3}f_i(\stackrel{}{q},t)\sigma _iv_{rel,i}\frac{p_\nu q_i^\nu }{EE_i^{}}\right\},$$
(16)
To obtain $`𝒩_{h.g.}(ϵ_0)`$ one needs only to integrate (16) over $`p_T`$:
$$𝒩_{h.g.}(ϵ_0)=𝑑p_Tg(p_T,ϵ_0)\mathrm{exp}\left\{_{t_0}^{t_{f.o.}}𝑑t\underset{i=1}{\overset{l}{}}\frac{d^3\stackrel{}{q}}{(2\pi )^3}f_i(\stackrel{}{q},t)\sigma _iv_{rel,i}\frac{p_\nu q_i^\nu }{EE_i^{}}\right\}.$$
(17)
Now we would like to take the final-size effects and the transverse expansion into account in our model. To do this directly, we would have to come back to the formula given by (13) and integrate it, instead of (16) , over $`p_T`$. But this would involve a five-dimensional integral (the three-dimensional integral over $`\stackrel{}{q}`$ simplifies to the one-dimensional one, in fact) instead of the three-dimensional integral of (17). Therefore, we need to simplify in some way the direct method just mentioned above. We shall define an average time of leaving the hadron medium by $`J/\mathrm{\Psi }`$’s with the velocity $`v`$ produced in an A-B collision at impact parameter $`b`$, $`t_{esc}(b,v)`$. Then, we will put this quantity, instead of $`t_{f.o.}`$, as the upper limit of the integral over $`t`$ in (17).
Let us consider an A-B collision at impact parameter $`b`$. Since we will compare final results with the latest data of NA50 which are for Pb-Pb collisions , we focus on the case of A=B here. So, for the collision at impact parameter $`b`$ we have the situation in the plane $`z=0`$ as presented in Fig. 2, where $`S_{eff}`$ means the area of the overlap of the colliding nuclei. We shall assume here, that the hadron medium, which appears in the space between the nuclei after they crossed each other also has the shape of $`S_{eff}`$ at $`t_0`$ in the plane $`z=0`$. And additionally, the transverse expansion will start in the form of the rarefaction wave moving inward $`S_{eff}`$ at $`t_0`$. Then, for a $`J/\mathrm{\Psi }`$ which is at $`\stackrel{}{r}S_{eff}`$ at the moment $`t_0`$ and has the velocity $`\stackrel{}{v}`$ we denote by $`t_{esc}`$ the moment of crossing the border of the hadron gas. It means that $`t_{esc}`$ is a solution of the equation $`\stackrel{}{d}+\stackrel{}{v}(tt_0)=R_Ac_s(tt_0)`$, where $`R_A=r_0A^{\frac{1}{3}}`$ ($`r_0=1.2fm`$) is the nucleus radius and $`\stackrel{}{d}=\stackrel{}{r}\stackrel{}{b}`$ for the angel between $`\stackrel{}{r}`$ and $`\stackrel{}{v}`$ such that the $`J/\mathrm{\Psi }`$ will cross this part of the edge of the area of the hadron gas which was created by the projectile and $`\stackrel{}{d}=\stackrel{}{r}`$ in the opposite. Having obtain $`t_{esc}`$, we average it over the angel between $`\stackrel{}{r}`$ and $`\stackrel{}{v}`$, i.e. we integrate $`t_{esc}`$ over this angel and divide by $`2\pi `$. Then we average the result over $`S_{eff}`$ with the weight given by
$$p_{J/\mathrm{\Psi }}(\stackrel{}{r})=\frac{T_A(\stackrel{}{r})T_B(\stackrel{}{r}\stackrel{}{b})}{T_{AB}(b)}$$
(18)
and we obtain $`t_{esc}(b,v)`$. So, the final expression for $`𝒩_{h.g.}(ϵ_0)`$ when the transverse expansion is taken into account reads
$$𝒩_{h.g.}(ϵ_0)=𝑑p_Tg(p_T,ϵ_0)\mathrm{exp}\left\{_{t_0}^{t_{esc}}𝑑t\underset{i=1}{\overset{l}{}}\frac{d^3\stackrel{}{q}}{(2\pi )^3}f_i(\stackrel{}{q},t)\sigma _iv_{rel,i}\frac{p_\nu q_i^\nu }{EE_i^{}}\right\}.$$
(19)
## IV The energy density in the CRR
To compare our theoretical estimations for $`J/\mathrm{\Psi }`$ survival factor with the experimental data we need to modify the latest so as they become $`ϵ_0`$-dependent instead of $`E_T`$-dependent ($`E_T`$ is the neutral transverse energy). To do this we will use the well-known Bjorken formula
$$ϵ_0=\frac{3E_T}{\mathrm{\Delta }\eta S_{eff}t_0},$$
(20)
where $`\mathrm{\Delta }\eta `$ is the pseudo-rapidity range. Using values of impact parameter $`b`$ given in we can calculate $`S_{eff}`$ for each measured $`E_T`$ bin.
In further considerations we will need the formula for the number of participating nucleons as a function of impact parameter $`b`$:
$$N_{part}(b)=_{S_{eff}}d^2\stackrel{}{s}\left\{T_A(\stackrel{}{s})+T_B(\stackrel{}{s}\stackrel{}{b})\right\}.$$
(21)
We found out that the ratio $`E_T`$/$`N_{part}`$ is almost constant and for NA50 data varies from 0.28 GeV to 0.31 GeV with the average value equal roughly to 0.3 GeV and the standard deviation equal to 0.0094 GeV. Therefore we assume that for Pb-Pb collisions of NA50 the following approximation is valid:
$$E_T0.3N_{part}.$$
(22)
Having put (22) into (20) we obtain $`ϵ_0`$ as a function of $`b`$
$$ϵ_0(b)=0.75\frac{N_{part}(b)}{S_{eff}(b)},$$
(23)
where we have also used the value $`\mathrm{\Delta }\eta =1.2`$ of NA50 . The above function is depicted in Fig. 3. The behaviour in low $`b`$ is the most interesting feature of $`ϵ_0(b)`$. We can see that for $`b7.9`$ there are two different values $`b_1`$ and $`b_2`$ such that $`ϵ_0(b_1)=ϵ_0(b_2)`$. Of course, this is the result of direct application of (22) which is some approximation in fact. But nevertheless this can suggest that there is a wide range of impact parameter $`b`$ for which the resulting $`ϵ_0`$ is almost constant. Using the dependence between $`E_T`$ and $`b`$ obtained by NA50 , we can state the similar for $`E_T`$, i.e. for $`E_T4050GeV`$ $`ϵ_0`$ changes weakly with $`E_T`$. Also the same is true for $`J/\mathrm{\Psi }`$ suppression what could suggest that the quantity of $`ϵ_0`$ reached is the main reason for the suppression.
## V Results
To evaluate formulae (17) and (19) we have to know $`T(t)`$, $`\mu _B(t)`$ and $`\mu _S(t)`$ and how to obtain these functions was explained in Sect. II. But to follow all that procedure we need initial values $`s_0`$ and $`n_B^0`$. To estimate initial baryon number density $`n_B^0`$ we can use experimental results for S-S or Au-Au collisions. In the first approximation we can assume that the baryon multiplicity per unit rapidity in the CRR is proportional to the number of participating nucleons. For a sulphur-sulphur collision we have $`dN_B/dy6`$ and 64 participating nucleons. For the central collision of lead nuclei we can estimate the number of participating nucleons at $`2A=416`$, so we have $`dN_B/dy39`$. Having taken the initial volume in the CRR equal to $`\pi R_A^21`$ fm, we arrive at $`n_B^00.25fm^3`$. This is some underestimation because the S-S collisions were at a beam energy of 200 GeV/nucleon, but Pb-Pb at 158 GeV/nucleon. From the Au-Au data extrapolation one can estimate $`n_B^00.65fm^3`$ . These values are for central collisions, and for the higher impact parameter (a more peripheral collision) the initial baryon number density should be much lower. In fact, if we apply the above-mentioned assumption, the initial baryon number density for a given collision at the impact parameter $`b`$ will be proportional to the number of participating nucleons divided by $`S_{eff}`$. Therefore, $`n_B^0(b)`$ will have exactly the same shape as $`ϵ_0(b)`$ presented in Fig. 3. As a result, only for the most peripheral collisions $`n_B^0`$ will be substantially below the value for the central one. So, to simplify numerical calculations we will keep $`n_B^0`$ constant over the all range of $`b`$ and additionally, to check the possible dependence on $`n_B^0`$, we will do our estimations for $`n_B^0`$ substantially lower, i.e. $`n_B^0=0.05fm^3`$. According to our approximation of $`n_B^0(b)`$, it would correspond to the most peripheral Pb-Pb collisions, $`b14fm`$.
Now, to find $`s_0`$, first we have to solve (3 \- 5) with respect to $`T`$, $`\mu _S`$ and $`\mu _B`$, where we put $`ϵ=ϵ_0`$, $`n_B=n_B^0`$ and $`n_S=0`$. Then, having put $`T`$, $`\mu _S`$ and $`\mu _B`$ into (6) we obtain $`s_0`$. Finally, expressing left sides of (4,6) by (7) and after then solving (4 \- 6) numerically we can obtain $`T`$, $`\mu _S`$ and $`\mu _B`$ as functions of time. In fact, evaluating formulae (17) and (19) we do the following: first, we calculate $`T=T(t)`$ which turns out to be very well approximated by the expression
$$T(t)T_0t^a$$
(24)
and then we put this approximation into (17) and (19) . And for $`\mu _S(t)`$ and $`\mu _B(t)`$ in $`f_i(\stackrel{}{q},t)`$ we put solutions of (4,5) where $`n_B`$ given by (7) and $`n_S=0`$ and $`T`$ is given by (24). But the exponent $`a`$ in (24) has proven not to be unique for the whole range of $`T_0`$ considered here. One gets different values of the initial energy density $`ϵ_0`$ for different values of the impact parameter $`b`$ and for different geometry of the collision process. So $`b`$ dependent $`a`$ gives also $`b`$ dependent freeze-out time $`t_{f.o}`$. The formula (23) is calculated for different values of $`b`$ and Eqs. (3 \- 5) are solved. We have evaluated the suppression factor up to $`ϵ_0=3.5GeV/fm^3`$. This gives e.g. the maximal possible $`T_0`$, $`T_{0,max}`$, equal to $`219.3MeV`$ ( for $`n_B^0=0.65fm^3`$), $`224MeV`$ (for $`n_B^0=0.25fm^3`$) or $`224.8MeV`$ (for $`n_B^0=0.05fm^3`$).
This procedure allows to evaluate $`J/\mathrm{\Psi }`$ survival factor given by (17) . Because of the lack of data, we shall assume only two types of the cross-section, the first, $`\sigma _b`$, for $`J/\mathrm{\Psi }`$-baryon scattering and the second, $`\sigma _m`$, for $`J/\mathrm{\Psi }`$-meson scattering. For $`\sigma _b`$ we put $`\sigma _b=\sigma _{J/\psi N}`$. As far as $`\sigma _m`$ is concerned, we assume that this cross-section is $`2/3`$ of the corresponding cross section for baryons, which is due to the quark counting. In the following, we will use values of $`J/\mathrm{\Psi }Nucleon`$ absorption cross-section $`\sigma _{J/\psi N}35mb`$ obtained from p-A data . At the beginning, to illustrate how the value of power $`a`$ influences $`J/\mathrm{\Psi }`$ suppression we present in Fig. 1 two results: the first for $`a=\frac{1}{3}`$ (which is the exact value for a free massless gas) and the second for $`a=\frac{1}{5.6}`$ (which is the approximate value for the hadron gas and $`T_0200MeV`$). We can see that the suppression improves more than twice for the highest $`ϵ_0`$ indeed.
To make our investigations more realistic we have to take into account that only about $`60\%`$ of $`J/\mathrm{\Psi }`$ measured are directly produced during collision. The rest is the result of $`\chi `$ ($`30\%`$) and $`\psi ^{}`$ ($`10\%`$) decay . Therefore the realistic $`J/\mathrm{\Psi }`$ survival factor should read
$$𝒩(ϵ_0)=0.6𝒩_{J/\psi }(ϵ_0)+0.3𝒩_\chi (ϵ_0)+0.1𝒩_\psi ^{}(ϵ_0),$$
(25)
where $`𝒩_{J/\psi }(ϵ_0)`$, $`𝒩_\chi (ϵ_0)`$ and $`𝒩_\psi ^{}(ϵ_0)`$ are given also by formulae (8 -19) but with $`p_T^2_{J/\mathrm{\Psi }}^{AB}(ϵ)=p_T^2_{J/\mathrm{\Psi }}^{AB}(ϵ),p_T^2_\chi ^{AB}(ϵ),p_T^2_\psi ^{}^{AB}(ϵ)`$, $`K_{J/\psi }=K_{J/\psi },K_\chi ,K_\psi ^{}`$, $`\sigma _{J/\psi N}=\sigma _{J/\psi N},\sigma _{\chi N},\sigma _{\psi ^{}N}`$ and $`M=M_{J/\psi },M_\chi ,M_\psi ^{}`$ respectively. The remaining problem is whether formula (9) is valid for $`\chi `$ and $`\psi ^{}`$. There are data for $`p_T^2_\psi ^{}^{PbPb}`$ and they shows that $`p_T^2_\psi ^{}^{PbPb}1.4p_T^2_{J/\mathrm{\Psi }}^{PbPb}`$. So, we assume that the above is also true for $`p_T^2_\psi ^{}^{AB}(ϵ)`$, i.e.
$$p_T^2_\psi ^{}^{AB}(ϵ)=1.4p_T^2_{J/\mathrm{\Psi }}^{AB}(ϵ)$$
(26)
with $`p_T^2_{J/\mathrm{\Psi }}^{AB}(ϵ)`$ given by (9) . For $`\chi `$ we believe that the inequality
$$p_T^2_{J/\mathrm{\Psi }}^{AB}p_T^2_\chi ^{AB}p_T^2_\psi ^{}^{AB}$$
(27)
should be valid and therefore assume that (26) is true also in this case. Anyway, the exact form of $`p_T^2_\chi ^{AB}(ϵ)`$ or $`p_T^2_\psi ^{}^{AB}(ϵ)`$ is not very important because we checked that the suppression depends on this form very weakly. First, we put $`K_{J/\psi }=0`$ and the resulting $`J/\mathrm{\Psi }`$ survival factor (for direct $`J/\mathrm{\Psi }`$’s) differs only a few percent for the highest $`ϵ_0`$ from the one calculated with formula (9) unchanged. Second, when we use expression (9) also for $`\chi `$ and $`\psi ^{}`$, the evaluated suppression factor is the same as that calculated with the use of (26), as far as plots are concerned.
To complete our estimations we need also values of cross-sections for $`\chi baryon`$ and $`\psi ^{}baryon`$ scatterings (we will still hold that $`\chi (\psi ^{})meson`$ cross-section is $`\frac{2}{3}`$ of $`\chi (\psi ^{})baryon`$ cross-section). Since $`J/\mathrm{\Psi }`$ is smaller than $`\chi `$ or $`\psi ^{}`$, $`\chi baryon`$ and $`\psi ^{}baryon`$ cross-sections should be greater than $`J/\mathrm{\Psi }baryon`$ one. For simplicity, we assume that all these cross-sections are equal. This means that we underestimate $`J/\mathrm{\Psi }`$ suppression, here. The final results of calculations of (17) are presented in Figs. 5-7 for various sets of parameters of our model (which are $`T_{f.o.},n_B^0,\sigma _b`$). We performed these calculations for two values of $`T_{f.o.}=100,140MeV`$ which agree fairly well with values deduced from hadron yields . For comparison, also the experimental data are shown in Figs. 5-7. The experimental survival factor is defined as
$$𝒩_{exp}=\frac{\frac{B_{\mu \mu }\sigma _{J/\psi }^{AB}}{\sigma _{DY}^{AB}}}{\frac{B_{\mu \mu }\sigma _{J/\psi }^{pp}}{\sigma _{DY}^{pp}}},$$
(28)
where $`\frac{B_{\mu \mu }\sigma _{J/\psi }^{AB(pp)}}{\sigma _{DY}^{AB(pp)}}`$ is the ratio of the $`J/\mathrm{\Psi }`$ to the Drell-Yan production cross-section in A-B(p-p) interactions times the branching ratio of the $`J/\mathrm{\Psi }`$ into a muon pair. The values of the ratio for p-p, S-U and Pb-Pb are taken from . Note that since the equality $`\sigma _{DY}^{AB}=\sigma _{DY}^{pp}AB`$ has been confirmed experimentally up to now , formula (28) reduces to
$$𝒩_{exp}=\frac{\sigma _{J/\psi }^{AB}}{AB\sigma _{DY}^{pp}},$$
(29)
which is also given as the experimental survival factor, for instance, in .
Coming back to examination of our first results presented in Figs. 5-7, we can see that the most of Pb-Pb data are below the region of suppression obtained for chosen values of parameters in our model. Generally, the theoretical curves decrease slower with increasing $`ϵ_0`$, whereas the data show rather abrupt fall just above $`ϵ_0=2GeV/fm^3`$. Note that the dependence on the initial baryon number density is substantial but for higher values of $`n_B^0`$, rather. The lower the initial baryon number density, the deeper the suppression. There are two reasons for such a behaviour: the first, for the higher baryon number density, there are less non-strange heavier mesons $`\rho `$, $`\omega `$ in the hadron gas of the same $`ϵ_0`$, but these particles create the most weighty fraction of scatters, for which reaction (1) have no threshold at all; the second, the freeze-out time $`t_{f.o.}`$ decreases with increasing $`n_B^0`$ for a given $`ϵ_0`$ in our model. For instance, for $`ϵ_0=3.5GeV/fm^3`$ and $`T_{f.o.}=140MeV`$ we have $`a=0.172,\mathrm{\hspace{0.33em}0.175},\mathrm{\hspace{0.33em}0.183}`$ and $`t_{f.o.}=15.7,\mathrm{\hspace{0.33em}14.7},\mathrm{\hspace{0.33em}11.6}fm`$ for $`n_B^0=0.05,\mathrm{\hspace{0.33em}0.25},\mathrm{\hspace{0.33em}0.65}fm^3`$, respectively. We can see also that the value $`\sigma _b=3mb`$ is too small to obtain results comparable with the data, so we will leave aside this value in further investigations.
Now we will include the finite-size effects into our model, i.e. we will take into account that the realistic hadron gas has a finite transverse size. This will be done in form of the rarefaction wave moving inward $`S_{eff}`$ with the sound velocity $`c_s`$. How to obtain this velocity has been mentioned in Sec.2 (see also ). With the finite-size effects included, the final expression for $`J/\mathrm{\Psi }`$ survival factor $`𝒩_{h.g.}(ϵ_0)`$ will be given by (19) . To make our investigations much more realistic we will also include the possible $`J/\mathrm{\Psi }`$ disintegration in nuclear matter, which should increase $`J/\mathrm{\Psi }`$ suppression by about $`10\%`$ . But to draw also S-U data in figures, instead of multiplying $`𝒩_{h.g.}`$ by $`𝒩_{n.m.}`$ given by (11), we divide $`𝒩_{exp}`$ by appropriate $`𝒩_{n.m.}`$, i.e. we define ”the experimental $`J/\mathrm{\Psi }`$ hadron gas survival factor” as
$$\stackrel{~}{𝒩}_{exp}=\mathrm{exp}\left\{\sigma _{J/\psi N}\rho _0L\right\}𝒩_{exp}.$$
(30)
and values of this factor are drawn in Figs. 9-10 as the experimental data.
At the beginning we will consider a uniform nuclear matter density
$$\rho _A(\stackrel{}{s},z)=\rho _A(\stackrel{}{r})=\{\begin{array}{cc}\rho _0=\left(\frac{4\pi }{3}r_0^3\right)^1,\stackrel{}{r}R_A\hfill & \\ \mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}0},\stackrel{}{r}>R_A\hfill & \end{array}.$$
(31)
The results of numerical estimations of (19) and (30) are depicted in Fig. 9 for two values of the initial baryon number density, $`n_B^0=0.25,\mathrm{\hspace{0.33em}0.65}fm^3`$. The curve for $`n_B^0=0.05fm^3`$ almost covers the curve for $`n_B^0=0.25fm^3`$ and the maximal difference between these curves do not exceed $`1.9\%`$ (which is the difference for the highest possible $`ϵ_0`$, i.e. $`ϵ_02GeV/fm^3`$), so for clearness of the figure we do not draw it. The two values of the speed of sound are the maximal values of this quantity possible in the range $`[T_{f.o.}=140MeV,T_{0,max}]`$ for the above-mentioned two cases of $`n_B^0`$. In fact, we have checked that the results almost do not depend on $`c_s`$ (allowed in the range) and the difference (seen only for the quantity $`ϵ_0`$) between $`J/\mathrm{\Psi }`$ survival factors for the maximal and the minimal values of $`c_s`$ in the range $`[T_{f.o.}=140MeV,T_{0,max}]`$ are less than $`0.6\%`$. Note that the theoretical curves in Figs. 9 (the same will happen in Figs. 9-10) are two-valued around $`ϵ_0=2GeV/fm^3`$. This is the result of our approximation of $`ϵ_0(b)`$ given by (23). This expression allows for two different values of $`b`$, which give the same $`ϵ_0`$ in some range of the impact parameter (see Sec.4).
It has turned out also that in the case of the transverse expansion, the results almost do not depend on the $`T_{f.o.}`$ (for $`T_{f.o.}[100,140]MeV`$). And the maximal difference (seen for the quantity $`ϵ_0`$) between curves for $`T_{f.o.}=100MeV`$ and $`T_{f.o.}=140MeV`$ do not exceed $`2.7\%`$ ($`n_B^0=0.05fm^3`$), $`2.6\%`$ ($`n_B^0=0.25fm^3`$). This is because the freeze-out time resulting from the transverse expansion, $`t_{f.o.,trans}=R_A/c_s`$ (if we assume a central collision and $`c_s`$ constant), is of the order of the freeze-out time resulting from the longitudinal expansion for $`T_{f.o.}=140MeV`$. Namely, for Pb and $`c_s=0.45`$ we have $`t_{f.o.,trans}15.8fm`$ which is very similar to values of $`t_{f.o.}`$ for $`T_{f.o.}=140MeV`$ given earlier. For $`T_{f.o.}=100MeV`$, $`t_{f.o.}=111.0,\mathrm{\hspace{0.33em}101.0},\mathrm{\hspace{0.33em}72.5}fm`$ for $`n_B^0=0.05,\mathrm{\hspace{0.33em}0.25},\mathrm{\hspace{0.33em}0.65}fm^3`$ respectively, so the hadron gas ceases because of the transverse expansion much earlier.
We repeated our estimations of formula (19) also for the more realistic nuclear matter density distribution, namely for the Woods-Saxon distribution with parameters taken from . The results are presented in Figs. 9-10. We can see that these curves fit the data a little bit better than those obtained within the assumption of the uniform nuclear matter density distribution (cf. Fig. 9 and 9). Generally, taking into account also the transverse expansion changes the final (theoretical) pattern of $`J/\mathrm{\Psi }`$ suppression qualitatively. First of all, the curves for the case including the transverse expansion are concave (what is more clearly seen for lower $`n_B^0`$) as the data suggest should be, in opposite to the case with the longitudinal expansion only, where curves are convex. But still, theoretical curves are not steep enough to cover the data area completely. We would like to note at this point that there is some ambiguity in calculation of $`E_T`$ and $`b`$ in NA50 experiment, since the range of $`ϵ_0`$ obtained from (20) with the use of $`E_T`$ and $`b`$ for 1995 and 1996 Pb-Pb runs are different. For 1995 run we have $`ϵ_02.62.9GeV/fm^3`$ and for 1996 run $`ϵ_02.02.1GeV/fm^3`$, both estimates are for $`E_T40GeV`$. Of course, where exactly the data points should be placed is crucial for the valuation of the correctness of the shape of theoretical curves.
When the additional disintegration in the nuclear matter is included, also the magnitude of $`J/\mathrm{\Psi }`$ suppression is comparable with the data, but rather for greater $`\sigma _b`$. But note that since we have one overall charmonium-baryon cross-section $`\sigma _b`$, our final results underestimate the suppression (for $`\chi ,\psi ^{}baryon`$ scattering the cross-section should be greater than for $`J/\mathrm{\Psi }`$).
As a final remark, we think that it is difficult to exclude $`J/\mathrm{\Psi }`$ scattering in the hot hadron gas entirely, as the reason for the observed $`J/\mathrm{\Psi }`$ suppression at this point. In our model the most crucial parameter is the charmonium-baryon inelastic cross-section and the final results depend on its value substantially. Therefore it is of the greatest importance to establish how this cross-section behaves in the hot hadron environment. Some work has been done into this direction , but results presented there differ from each other and are based on different models.
## Acknowledgements
We would like to thank Dr K.Redlich for very helpful discussions. |
warning/0002/gr-qc0002082.html | ar5iv | text | # 1 Evolution of the scale factor in case A₂=-1. Initial condition is follows: R(0)=0,{dR/dt|}ₜ₌₀=1.
Effect of latent mass in inhomogeneous cosmological model with perfect fluid and self-acting scalar field.
V.M. Zhuravlev<sup>1</sup>, D.A. Kornilov<sup>2</sup>
Ulyanovsk State University
Institute for Theoretical Physics
<sup>1</sup> zhuravl@sv.uven.ru
<sup>2</sup> kda@sv.uven.ru
Abstract. The inhomogeneous cosmological model with matter in the form self-acting scalar field and perfect fluid is considered. On the basis of exact solutions is considered the evolution of density distribution of a matter in space on a background cosmological expansion by the Universe. Is shown, the first, in such model the equation of a matter state is variable in time and is closely connected to character cosmological expansions. Secondly, it is shown with point of view of the observer the Universe looks as space flat, but with effect of latent mass. This effect consists in that the mass of a perfect fluid by dynamic measurements surpasses the own mass perfect fluid that is explained by presence scalar field.
PACS-94: 04.20.-q
The usual approach to study cosmological models consists from two stages. The first stage consists in construction of global homogeneous and isotropic model of the Universe of FRW-type with a scalar field and (or) matter in the form perfect fluid. At the second stage it is studied evolution perturbation of matter density on a background of the homogeneous metrics constructed at the first stage. Thus the inhomogeneity of the Universe is taken into account only in the first order of the perturbation theory, that is connected with linearization of the dynamics equations. However, finding out structure of the Universe on the large scales (voids, strings) show that the process of their formation can be essentially non-linear and for simulation of their occurrence it is necessary to use some other approaches. In the present work the attempt is undertaken to investigate properties of the self- coordinated model of a special scalar field and perfect fluid inducing a gravitational field on the basis of the exact solutions of the Einstein equations, describing Universe with cosmological evolution. In this case the distribution of matter in the Universe are described exact, instead of by the approximate solutions of the equations. Therefore in such model it is possible to study all effects of interaction of three fields in an exact form in formation various spatial structures of a matter in space.
The most suitable material for construction this models is well-known in the General Relativity theory the Majumdar-Papapetrou metrics originally used for construction self-gravitating electrostatic or magnitoststic fields.
Basis of suggested in the present work models is the generalization of the Majumdar-Papapetrou metrics to describe not static effects connected with cosmological expansion. Thus an electrostatic or magnitoststic field is replaced by some generalized scalar field. Nevertheless this scalar field has a similar energy-momentum tensor to electrostatic field distinguished from them only by presence of its self-action. Self-action of this field describes by potential function which one in many respects determines cosmological properties of model. The second kind of material field in this model is the perfect fluid. In result the researched model represents inhomogeneous cosmological model with self-acting scalar field and matter with state equation generally varies in due course. Last fact is important from the point of view modern representations about existence of various epochs in evolution of the Universe during the substance filling the Universe had various properties. At an inflationary stage it is a matter close to quasi-vacuum state, i.e. it prevails of a field component, in the subsequent epoch it is the isotropic radiation ($`p=\epsilon /3`$) and now it is a matter close to a dust ($`p0`$). In the present work the attempt is done(made) to analyse such inhomogeneous models from the point of view of evolution including all these stages of development by the Universe.
1. General view of the metrics researched in the present work is following
$`ds^2=e^{A(x,y,z)+b(t)}dt^2e^{A(x,y,z)+a(t)}\left(dx^2+dy^2+dz^2\right),`$ (1)
$`g_{ik}=\mathrm{diag}\{e^{A+b},e^{A+a},e^{A+a},e^{A+a}\},i,k=0,1,2,3,`$
where $`A=A(x,y,z)`$ \- function of coordinates $`x^1=x,x^2=y,x^3=z`$ and not dependent from $`x^0=t`$, and $`a(t),b(t)`$ \- some functions of time. The metric of such form with $`a(t)=0,b(t)=0`$ is known in the classical theory of GR as the Majumdar-Papapetrou metrics. In a case $`a(t)0,b(t)0`$ these not static metrics also describe some global cosmological dynamics of the Universe with local inhomogeneity of the space it was connected with function $`A`$. There is one more important underclass of the metrics of this type having the following form
$`ds^2=e^{A(x,y,t)+a(z)}\left(dt^2dx^2dy^2\right)e^{A(x,y,t)+b(z)}dz^2.`$ (2)
It describes non-static gravitational processes in space with coordinate $`z`$. Dependence from $`z`$ of functions $`a,b`$ in this case is connected to some inhomogeneous and anisotropy space lengthways the select axis $`z`$.
Let’s consider the matter inducing metric property space - time of the kind (1) the mix from perfect fluid with energy-momentum tensor (TEM)
$$T_0^{(m)0}=\epsilon (x,y,z,t),T_1^{(m)1}=T_2^{(m)2}=T_3^{(m)3}=p(x,y,z,t),T_k^{(m)i}=0,ik$$
(3)
($`\epsilon `$ and $`p`$ \- density of energy and pressure of a fluid) and scalar field with energy-momentum tensor
$$T_{ik}^{(\varphi )}=_i\varphi _k\varphi +\frac{1}{2}g_{ik}g^{lm}_l\varphi _m\varphi +g_{ik}V(\varphi ,t)$$
(4)
( $`V(\varphi ,t)`$ \- potential of self-action of a scalar field). As follows from (4) TEM of a scalar field considered in this work differs from TEM of usual scalar field, for example Higg’s scalar field, by opposites sign, but coincides with sign of TEM of an electrostatic field with potential $`\varphi `$. However, equation (4) differs from TEM of an electrostatic field by presence of self-action potential.
The Einstein equations
$$G_{ik}=\ae T_{ik}$$
(5)
( $`\ae =8\pi G/c^4`$ \- Einstein’s gravitational constant) for the metrics (1) with TEM
$$T_{ik}=T_{ik}^{(\varphi )}+g_{ij}T_k^{(m)j}$$
are reduced to simple set from two equations, which is possible to write down in the following form
$`p(x,y,z,t)=V(\varphi )+g(t)e^{bA},`$ (6)
$`\epsilon (x,y,z,t)=c^2\rho (x,y,z,t)={\displaystyle \frac{1}{\ae }}e^{Aa}\mathrm{\Delta }AV(\varphi )+{\displaystyle \frac{\dot{a}^2}{\ae }}e^{bA},`$ (7)
$$g(t)=\frac{1}{\ae }\left[\frac{1}{2}\dot{a}(\dot{b}\dot{a})\ddot{a}\right],$$
and one statement
$$\varphi (x,y,z,t)=\frac{a(t)A(x,y,z)}{\sqrt{2\ae }},$$
(8)
identifying a scalar field with characteristic of the metrics. To these equations it is necessary obviously to add the equation arising the ambassador variations of a Lagrangian density of a matter by $`\varphi `$. This equation has the following form
$$\mathrm{\Delta }A=\sqrt{2\ae }e^{A+a}\frac{V(\varphi )}{\varphi }+\left(\frac{\dot{b}3\dot{a}}{2}\dot{a}\ddot{a}\right)e^{ab2A}$$
(9)
Further, for function $`A`$ the equation (9) should not on the right contain obvious dependence from coordinate $`t`$, as function $`A`$ is obvious from $`t`$ does not depend. It imposes specific conditions on potential $`V(\varphi ,t)`$. By the elementary kind $`V(\varphi ,t)`$ it is possible to satisfy to condition of independence of the right part in (9) from $`t`$, is potential of a kind
$$V(\varphi )=V_0\mathrm{exp}\{\sqrt{2\ae }\varphi \}.$$
(10)
In this case we have
$`\mathrm{\Delta }A=\sigma e^{2A},`$ (13)
$`p=g(t)e^{bA}+V(\varphi )=p_0(t)e^A,`$
$`\epsilon =c^2\rho =g(t)e^{bA}V(\varphi )+\sqrt{{\displaystyle \frac{2}{\ae }}}{\displaystyle \frac{V}{\varphi }}=p_0(t)e^A,`$
where
$`p_0(t)=g(t)e^b+V_0e^a,`$
$`\sigma =2\ae V_0e^{2a}+\left({\displaystyle \frac{\dot{b}3\dot{a}}{2}}\dot{a}\ddot{a}\right)e^{ab}=\mathrm{const}.`$
Last statement is equation connecting $`a`$ and $`b`$, at any meanings of parameters $`\sigma `$ and $`V_0`$. From here for (10) we come to the extreme rigid state matter equation
$$p=\epsilon ,$$
(14)
or to absence of a matter: at $`p=0`$ is automatically received $`\epsilon =0`$.
That the state matter equation would have more general view, for example, $`p=\gamma (t)\epsilon `$, it is necessary to require performance of the following equation for $`V(\varphi ,t)`$:
$$V(\varphi ,t)+q(t)\mathrm{exp}\{\sqrt{2\ae }\varphi \}=\gamma (t)\left(V(\varphi ,t)+\sqrt{\frac{2}{\ae }}\frac{V}{\varphi }+q(t)\mathrm{exp}\{\sqrt{2\ae }\varphi \}\right).$$
Here
$$q(t)=g(t)e^{b(t)a(t)}.$$
The general solution of this equation rather $`V(\varphi ,t)`$ at parametrical dependence from $`t`$ has the following form
$$V(\varphi ,t)=V_1(t)\mathrm{exp}\left\{\frac{\gamma +1}{2\gamma }\sqrt{2\ae }\varphi \right\}q(t)\mathrm{exp}\{\sqrt{2\ae }\varphi \},$$
(15)
where $`V_1(t)`$ \- any function $`t`$. The parameter $`\gamma `$ can be thus function of time, and can and to not be. However in all cases, when $`\gamma 1`$ and $`\gamma 0`$ the self-action potential will be obvious function of time. For example, for the case, when the substance represents by itself isotropic radiation with $`\gamma =1/3`$ potential looks like
$$V(\varphi ,t)=V_1(t)\mathrm{exp}\{2\sqrt{2\ae }\varphi \}q(t)\mathrm{exp}\{\sqrt{2\ae }\varphi \}.$$
(16)
In general case equation for $`A(x,y,z)`$ will look like
$$\mathrm{\Delta }A=A_1e^{(1+3\gamma )A/(2\gamma )}+A_2e^{2A}.$$
(17)
Beacause the right member of the equation (17) did not contain a time $`t`$ dependence, it is necessary that the functions $`a,b,V_1`$ satisfied to following requirements:
$$A_1=\mathrm{const}=4\ae V_1(t)e^{3a},$$
(18)
$$A_2=\mathrm{const}=\left[\frac{\dot{a}(\dot{b}+\dot{a})}{2}\ddot{a}\right]e^{ab}.$$
(19)
At once from (18) it is possible to obtain a kind of function $`V_1(t)`$:
$$V_1(t)=\frac{A_1\gamma }{\ae (1+\gamma )}e^{(1+3\gamma )a/(2\gamma )}.$$
(20)
Equation (19) contains two unknown functions, one of them $`(b)`$ remains uncertain and connected with selection of time variable. For simplicity let’s assume $`b(t)=0`$. It is means that variable $`t`$ is the phisical time. In this case differential equation
$$\ddot{a}\frac{\dot{a}^2}{2}=A_2e^a,$$
defines form of function $`a(t)`$, and it has the solution in an explicit form in only elliptic functions of the first kind .
Last equation determines evolution of the scale factor $`R(t)=\mathrm{exp}\{a(t)/2\}`$ too. In figure 1 the results of a numerical analysis of the differential equation
$$\frac{d^2R}{dt^2}=\frac{A_2}{2}R^3,$$
(21)
that defines evolution of the scale factor.
From the equation (21) it is visible that the changes of evolution of the scale factor are determined by the sign of the constant $`A_2`$. In particular, in a case $`A_2<0`$, the second derivative on time from the scale factor is positive (positive acceleration) and the expansion of the Universe will be accelerated. For more best estimate of velocity of expansion of the Universe in a fig. 1 the diagrams power and exponential functions are given. It is visible, that the expansion at a small times is a very good aproximated by a linear function, then the very fast growth of the scale factor follows.
2. To find out properties of a field $`\varphi `$ and its role in represented theory it is important to study dynamics of a trial particle in a gravitational field described by the metrics (1). Let’s assume, that the field $`\varphi `$ does not interact directly with a usual matter. Then at absence of direct interaction trial particle with a field $`\varphi `$ the equations of its dynamics will look like Geodetic
$$\frac{d}{dt}u^i+\mathrm{\Gamma }_{kj}^iu^ku^j=0,\frac{dx^k}{dt}=u^k,i,j,k=0,1,2,3,$$
(22)
where $`u^k`$ \- 4-speed of a particle and $`G_{kj}^i`$ is the Christoffel simbols. For the metrics (1) 4-speed of a particle are normed by a condition
$$u^iu_i=g_{ik}u^iu^k=e^A\left(e^{2A}(u^0)^2(u^1)^2(u^2)^2(u^3)^2\right)=c^2.$$
(23)
where $`c`$ \- speed of light.
Substituting in (22) Christoffel symbols for the metrics (1) we receive
$`{\displaystyle \frac{d}{dt}}u^\alpha ={\displaystyle \frac{1}{2}}\left(e^{2A}(u^0)^2+(u^1)^2+(u^2)^2+(u^3)^2\right)A_{,\alpha }+u^\alpha (u^\beta A_{,\beta }),`$ (24)
$`\alpha ,\beta =0,1,2;`$
$`{\displaystyle \frac{d}{dt}}u^0=u^0(u^\beta A_{,\beta }).`$ (25)
Having copied these equations for covariant components of 4-speed of a particle using statement
$$u^0=e^Au_0,u^1=e^Au_1,u^2=e^Au_2,u^3=e^Au_3,$$
we receive
$`{\displaystyle \frac{d}{dt}}u_0=0,`$ (26)
$`{\displaystyle \frac{d}{dt}}u_\alpha ={\displaystyle \frac{1}{2}}\left(2(u_0)^2c^2e^A\right)A_{,\alpha },\alpha =0,1,2.`$ (27)
Thus covariant component $`u_0`$ of 4-speeds of a particle is constant and last three equations are equivalent to equations of movement of a particle in a Newtonian field of gravitation:
$$\frac{d}{dt}u_\alpha =\mathrm{\Psi }_{,\alpha }\alpha =0,1,2,$$
(28)
with gravitation potential
$$\mathrm{\Psi }=(u_0)^2A\frac{c^2}{2}e^A.$$
(29)
Actually first member in this expression proportional to a square of kinetic energy of a trial particle is the item connected with potential forces of inertia and only second can be interpreted as Newtonian potential of a field of gravitation. For interpretation of movement the particles in this case are necessary for considering the equations (24). This equations is possible to write down as
$$\frac{d}{dt}u^\alpha =\frac{\mathrm{\Phi }}{x^\alpha }+F^\alpha ,$$
(30)
where $`\mathrm{\Phi }`$ now true potential of a field of gravitation:
$$\mathrm{\Phi }=\frac{c^2}{2}e^A+\mathrm{\Phi }_0,$$
(31)
$`\mathrm{\Phi }_0`$-any constant, and $`𝐅`$ is the gyrotropic force of inertia arising at the expense of local rotation of reference system. It will easily be convinced as $`𝐅`$ looks like
$$𝐅=[𝐮\times [𝐮\times A]],$$
where on the right there is a double vector product on flat contravariant card of co-ordinates. As the force $`𝐅`$ is gyrotropic, it does not make work.
Repeating the calculations have been carried out above for the metrics (1) it is possible to receive similar to (30) equations for contravariant and covariant component of 4-speeds in the metrics (2).
3. We shall consider now the distribution of a matter corresponds to potential fields of gravitation (31). Using (7),(10) we come to the following equation:
$$\mathrm{\Delta }\mathrm{\Phi }\frac{c^2}{2}e^A\left((A)^2+\mathrm{\Delta }A\right)$$
(32)
Last member in the right part, square-law on a gradient of a field, will be equal to energy ensity $`\epsilon _\varphi `$ of a scalar field $`\varphi `$. Using (7) the equation (32) it is possible to write down in a standard form
$$\mathrm{\Delta }\mathrm{\Phi }=4\pi G\varrho ,$$
where
$$\varrho =\rho _d+\frac{c^2}{2}V(\varphi )+c^2(\varphi )^2e^{\sqrt{2\ae }\varphi }$$
is complete density of all kinds of a matter in space: $`\rho _d`$ \- density usual matter, and the other items represent energy density of a field $`\varphi `$ at the expense of self-action of a field $`\varphi `$.
Let’s notice this result describe the presence of effect of latent mass. Scalar field in this model not detected by not dynamic measurements, as is actually connected by virtue of equality (8) with the metrics of space - time and itself does not interract with usual matter. But at dynamic measurements will be present additional mass. For any volume $`V`$ of space the additional mass $`M_l`$ is equal to following
$$M_l=\underset{\mathrm{\Omega }}{}\left[\frac{c^2}{2}V(\varphi )+c^2(\varphi )^2e^{\sqrt{2\ae }\varphi }\right]𝑑\mathrm{\Omega }.$$
This formula show that in considarating model the effect of latent mass exists. The value of latent mass define by function of $`A(x,y,z)`$ and dependent from cosmological variation of scale factor that equail to $`R(t)=\mathrm{exp}a(t)/2`$.
Acknowledgements. Work is supported by Russian Fund of Basic Researches (grant 01-98-18040). |
warning/0002/astro-ph0002125.html | ar5iv | text | # The Lyman Continuum Polarization Rise in the QSO PG 1222+2281footnote 11footnote 1Accepted for publication in Publications of the Astronomical Society of the Pacific, 2000 May
## 1 INTRODUCTION
Spectropolarimetric observations with the Hubble Space Telescope (HST) have revealed an unexpected rise in linear polarization in several QSOs (see review by Koratkar & Blaes 1999). These are radio quiet “candidate Lyman edge QSOs”, in which the continuum flux drops rather rapidly at rest wavelengths $`\lambda <1000`$ Å (Antonucci, Kinney, & Ford 1989; Koratkar, Kinney & Bohlin 1992). Models of accretion disk atmospheres predicted a reduced polarization in the Lyman continuum because of a diminished contribution of electron scattering to the opacity (Laor, Netzer, & Piran 1990). With this motivation, Impey et al. (1995) and Koratkar et al. (1995) used the Faint Object Spectrograph (FOS) on HST to obtain ultraviolet spectropolarimetry of several QSOs with redshifts sufficient to bring the Lyman continuum within the observed wavelength band. The surprising result, in several cases, was a rapid rise in polarization in the Lyman continuum. From values $`1`$ percent in the optical and near ultraviolet, the observed polarization rises around rest wavelength 750 Å to values $`5`$ percent in several objects, and to $`20`$ percent in PG 1630+377.
This phenomenon has inspired several attemps at explanation. Blaes and Agol (1996) found that, for effective temperatures $`T_{eff}25,000`$ K and low effective gravities, a polarization rise of up to $`5`$ percent at about the observed wavelength could occur naturally in QSO disk atmospheres. This results from the interplay of electron scattering, bound-free opacity, and the temperature gradient in the atmosphere. However, Shields, Wobus, and Husfeld (1998, hereinafter SWH) showed that the effects of the relativistic transfer function destroy the agreement between this model and observation. Beloborodov and Poutanen (1999) suggested a model involving Compton scattering in a corona or wind, but this model appears to have trouble giving the rapid rise in polarized flux observed in PG 1630+377 (Blaes & Shields 1999). Lee and Blandford (1997) discussed the possible role of scattering by resonance lines of heavy elements (see Section 5).
SWH showed that, if the polarization is assumed to rise sharply at $`\lambda 912`$ in the rest frame of the orbiting gas, then relativistic effects would naturally produce the wavelength dependence of the observed polarization. This may offer a way of measuring the black hole spin, but the physical mechanism for the polarization rise remains unknown.
PG 1222+228 is a $`\mathrm{B}15.5`$ radio quiet QSO (Schmidt & Green 1983) whose polarization rise at $`\lambda 750`$ Å coincides with a sharp drop in flux (Fig. 1, 2). Impey et al. (1995) noted this and attributed it to a coincidental Lyman limit system (LLS), corresponding to an identified absorption line system at z = 1.486. However, the coincidence of a broad absorption feature with a polarization rise also is observed for broad absorption line (BAL) QSOs. In these objects, outflowing gas at velocities $`10^4\mathrm{km}\mathrm{s}^1`$ produces blueshifted absorption troughs, typically seen in the resonance lines of H I, C IV, N V, O VI, Si IV, and sometimes Mg II (Weymann et al. 1991; Arav, Shlosman, & Weymann 1997). Spectropolarimetric observations (e.g., Ogle 1997; Schmidt & Hines 1999; Ogle et al. 1999) often show a rise in polarization in the troughs, reaching values as high as $`8`$ to 10 percent from $`1`$ percent at unabsorbed wavelengths. This is explained in terms of scattering of some of the continuum by an extended region that is not covered by the BAL flow (Hines & Wills 1995; Goodrich & Miller 1995; Cohen et al. 1995). This pattern resembles the polarization rise and flux drop in PG 1222+228.
This paper addresses two questions: (1) Does the polarization rise in PG 1222+228 result from an intrinsic absorber, analagous to the situation in the BAL QSOs? (2) If the drop in flux in PG 1222+228 is an intervening LLS, what are the consequences of correcting the observed, polarized continuum for this absorption?
## 2 INTRINSIC ABSORPTION IN PG 1222+228?
We first consider the possibility that the flux drop and coincidental polarization rise in PG 1222+228 results from some kind of intrinsic absorption. Two possibilities, considered below, are that it is a BAL outflow, or that it is an unusual, intrinsic LLS.
In either case, one issue is the behavior of the polarized flux as a function of wavelength. As discussed above, if the polarization rise results from the selective absorption of the directly viewed continuum but not the scattered continuum, one might expect a smaller drop (but generally not a rise) in the polarized flux, $`I_p=pI_\lambda `$. In order to examine this, we have rebinned the data of Impey et al. (1995), kindly made available in reduced form by C. Impey and C. Petri (1999). These data consist of a spectrum with the G190H grating covering $`\lambda \lambda 1575`$ to 2320 at 0.37 Å per pixel, and a spectrum with the G270H grating covering the range 2224 to 3295 at 0.52 Å per pixel. The G190H data shortward of $`\lambda 1994`$ have low signal-to-noise and were not presented by Impey et al. (1995). We used seven wavelength bins (in Å): (1)1994–2224, (2) 2224–2287, (3) 2287–2319, (4) 2319–2492, (5) 2492–2761, (6) 2761–3029, (7) 3029–3295. Bins 2 and 3 involve an average of the two overlapping spectra, and bin 3 is a narrow bin containing the flux drop at $`\lambda 2300`$. The resulting values of $`I_\lambda `$ , $`qQ/I`$ and $`uU/I`$ are tabulated in Table 1, along with the polarization, $`p=\sqrt{q^2+u^2}`$, and its position angle, $`\theta `$. For the polarized flux, we use the rotated Stokes flux $`Q^{}`$ and polarization $`q^{}Q^{}/I`$ (cf. Koratkar etal 1995), referred to a position angle of 168 degrees. This is based on the mean polarization position angle of the HST data, which is in reasonable agreement with optical observations (Stockman et al. 1984; Webb et al. 1993). (The use of $`q^{}`$ is appropriate if the position angle of the polarization is constant with wavelength. Table 1 supports this and also shows that the polarization $`p`$ is reasonably consistent with $`q^{}`$.) These quantities are plotted in Figure 2. The shortest wavelength bin has larger polarized flux than the longer wavelength bins. At face value, this would weigh against an intrinsic absorber model for PG 1222+228; but it involves a single wavelength bin with substantial error bars. Therefore, we consider other aspects of the two outflow models.
### 2.1 BAL Absorption
The rest wavelength of the onset of the absorption feature in PG 1222+228 is $`750`$ Å. Some BAL QSOs show absorption by Ne VIII $`\lambda `$775 (e.g., Arav et al. 1999; Telfer et al. 1999). This might be a candidate for the feature in question, in as much as BALs often set in at a wavelength somewhat blueshifted from the emission-line redshift. However, Ne VIII normally is accompanied by absorption in O VI $`\lambda 1035`$, N V $`\lambda 1240`$, and C IV $`\lambda 1550`$. There is no indication of broad C IV or Mg II absorption in the spectrum of PG 1222+228 (Sargent, Steidel, & Boksenberg 1988; Steidel & Sargent 1992). The HST spectrum shows a shallow trough at $`\lambda 3000`$ to $`\lambda 3070`$ that could be a weak O VI feature; but this may simply be a cluster of lines, including several strong L$`\alpha `$ lines indentified by Impey et al. (1996). Photoionization models by Hamann (1997) indicate a range of ionization parameters for which the fractional abundance of Ne<sup>+7</sup> exceeds that of O$`^{+5}.`$ However, given the normal ratio of oxygen to neon abundances, the O VI feature would likely be strong in a situation giving strong Ne VIII.
The flux drop at $`\lambda 750`$ does not recover, with decreasing wavelength, in a way suggestive of a BAL (Figure 1). The spectrum has not fully recovered by rest wavelength 650 Å, corresponding to an an outflow velocity of more than 40,000 $`\mathrm{km}\mathrm{s}^1`$ if attributed to Ne VIII. Some moderately narrow “mini-BAL” features have been observed at such high velocities (Hamann et al. 1997), but true BAL troughs rarely reach such velocities. The same can be said in connection with the possibility that the $`\lambda 750`$ feature corresponds to a blend of features including N III, N IV, O IV, S VI, and Ne VIII seen in BAL QSO spectra at this wavelength (e.g., Arav et al. 1999). Moreover, given the range of ionization stages contributing to this blend, C IV and Si IV absorption would likely accompany it.
Recent work has shown that BAL QSOs systematically have weak soft X-ray emission. BAL QSOs have optical to X-ray slopes $`\alpha _{ox}2.0`$, whereas nonBAL QSOs tend to have $`\alpha _{ox}`$ in the range -1.3 to -1.8 (Brandt, Laor, and Wills 1999). Here, $`\alpha _{ox}`$ is defined by $`F_x/F_o=(\nu _x/\nu _o)^{\alpha _{ox}}`$, where $`F_x`$ and $`F_o`$ are the flux densities ($`F_\nu `$) at 2 keV and 3000 Å, respectively. ROSAT pointing observations give a flux of $`6\times 10^{14}\mathrm{erg}\mathrm{s}^1\mathrm{cm}^2`$ for PG 1222+228 at a significance level of $`3.3\sigma `$ (Mushotzky 1999). If we assume a “normal” power-law slope of $`\alpha =1.6`$ over the 0.2 to 2 keV ROSAT band (Brandt et al. 1999), we find $`F_x=1.5\times 10^{31}\mathrm{erg}\mathrm{s}^1\mathrm{cm}^2\mathrm{Hz}^1`$ at 2 keV rest energy. Optical spectrophotometry (Wampler & Ponz 1985; Bechtold et al. 1984) implies $`F_o1.1\times 10^{26}\mathrm{erg}\mathrm{s}^1\mathrm{cm}^2\mathrm{Hz}^1`$ at rest wavelength 3000 Å. (These are observed fluxes at the wavelength corresponding to the indicated rest wavelength.) From this, we find $`\alpha _{ox}=1.8`$ for PG 1222+228. This result is uncertain because of the marginal X-ray detection and the possibility of variability, but at face value it is more consistent with a non-BAL than a BAL QSO. (Wilkes et al. 1994 quote an uncertain value $`\alpha _{ox}1.8`$ from Einstein data.)
We conclude that the flux drop at $`\lambda 750`$ in PG 1222+228 is unlikely to be a BAL feature.
### 2.2 An Instrinsic Lyman Limit System?
Impey et al. (1995) suggested that the flux drop at $`\lambda 750`$ was an intervening LLS. We show below that Lyman limit absorption does indeed give a good fit to the spectrum. This fit, however, leaves open the question of the location of the absorbing gas. Because of the coincidence with the polarization rise, we consider here the possibility that the feature is an intrinsic LLS, associated with a high velocity outflow from the QSO. Impey et al. (1996) identify L$`\alpha `$ absorption lines at z = 1.4857, 1.5238, 1.5272, and 1.5650 that might be associated with the $`\lambda 750`$ feature, if it is taken to be a LLS. A redshift of 1.486 corresponds to an outflow velocity $`60,000\mathrm{km}\mathrm{s}^1.`$ As noted above, this is not unprecedented for a QSO outflow producing absorption lines. However, the lines associated with the redshift systems in question in PG 1222+228 are narrow, and narrow lines with relative velocities greater than 5000 $`\mathrm{km}\mathrm{s}^1`$ usually are assumed to be intervening. For a Lyman edge optical depth of unity, the measured equivalent widths of the L$`\alpha `$ lines are consistent with a Doppler parameter $`b30\mathrm{km}\mathrm{s}^1`$ (see below), normal for an intervening LLS. In contrast, the mini-BALs observed at such high velocities have widths of order $`1000\mathrm{km}\mathrm{s}^1`$ (Hamann et al. 1997). Could the feature in PG 1222+228 nevertheless be caused by ejected material with a high outflow velocity and a small velocity dispersion?
We are not aware of any other case in which a high velocity LLS with narrow lines has been shown to be intrinsic. However, the existence of many intrinsic, narrow, high velocity absorption systems in QSOs has been proposed by Richards et al. (1999). These authors compare the incidence of absorption-line systems, per unit relative outflow velocity, for highly luminous QSOs with that for less luminous ones. In the velocity range 5000 to 75,000 $`\mathrm{km}\mathrm{s}^1`$, they find that absorption systems have a substantially larger frequency in luminous systems. Since intervening systems should have no dependence on QSO luminosity (assuming that discovery systematics are accounted for), Richards et al. conclude that at least the excess number of systems in the high luminosity QSOs are intrinsic.
A peculiar absorption line ratio in the z = 1.94 system in PG 1222+228 has been noted by Ganguly et al. (1998). These authors present high resolution spectra that show two narrow components, separated by $`17\mathrm{km}\mathrm{s}^1`$, with very different strengths of the Al II and Al III absorption lines. Photoionization models indicate that the component with strong aluminum lines must have an anomalously high abundance of aluminum. This is reminiscent of claims of unusual abundances in BAL QSOs, including an excess of aluminum (e.g., Turnshek et al. 1996; Junkkarinen et al. 1997; Shields 1997). The reality of these abundance anomalies is in doubt, because of the effects of partial covering of the continuum source (Arav 1997). However, the basic observation of anomalously strong Al lines may be a possible parallel between PG 1222+228 and the BAL QSOs, where outflowing gas is clearly present. If this is a hint that the narrow, z = 1.94 system may be intrinsic, perhaps it adds plausibility to the idea that the z = 1.486 system (or its neighbors) may also be intrinsic.
What might be the geometry of an intrinsic LLS in PG 1222+228? In order to explain the polarization rise, the absorber would have to intercept the line of sight to the continuum source but not the scattering source. The latter is often attributed to a wind driven off the inner edge of a “dust torus” (Krolik & Begelman 1986). The location of this may be related to the dust sublimation radius, $`0.2L_{46}^{1/2}`$ pc, where $`L_{46}`$ is the bolometric luminosity of the central source in units $`10^{46}\mathrm{erg}\mathrm{s}^1`$ (Laor & Draine 1993). The absorbing material, at a velocity of $`0.2c,`$ would be at a smaller radius in order not to cover the scattering source. Its high velocity suggests an origin at a small radius where the escape velocity is of order the observed outflow velocity. The escape velocity from a central mass of $`M`$ is 0.2$`c`$ at a radius $`10^{15.6}M_9`$ cm, where $`M_9M/10^9\mathrm{M}_{}`$. Let us assume the absorbing material is at a radius of $`10^{18}`$ cm, large enough to obscure the ultraviolet emitting part of an accretion disk but not the scattering source. At the observed speed, the crossing time would be a few years or less. Thus, the material should change radius substantially in the seven years between the Sargent et al. (1988) and the Impey et al. (1996) observations, and one might expect some change of velocity. These authors, however, quote velocities for the $`z=1.486`$ and 1.524 systems that agree within $`\mathrm{\Delta }z0.001`$. This corresponds to a change in outflow velocity of less than $`100\mathrm{km}\mathrm{s}^1`$. Such precise stability seems difficult to achieve. (Note, however, the stability of narrow features within the BAL profiles of some QSOs \[cf., Weymann 1997\]). A further problem involves the narrowness of the absorption lines. The radius of the ultraviolet emitting part of the disk would be at least $`10`$ gravitational radii, or about $`10^{15.1}M_9`$ cm. If the absorber is at a radius $`10^{18}`$ cm, then the line of sight through the absorber to different parts of the continuum source would likely give noticeably different projected flow velocities. The observed linewidths (see below) are $`30\mathrm{km}\mathrm{s}^1`$, less than one thousandth of the outflow velocity.
These difficulties with the intrinsic absorber model for PG 1222+228 encourage us to examine the straightforward idea of an intervening LLS.
## 3 INTERVENING ABSORPTION
Impey et al. (1995, 1996) attributed the flux drop at observed wavelength $`\lambda 2300`$ to a LLS associated with the $`z=1.4857`$ absorption line system. They identify absorption lines of H I, C III, N I, N II, Si II, and Si III. They also identify systems at 1.5238, 1.5272, and 1.5650, which have multiple Lyman lines and, for $`z=1.5238`$, C III. Sargent, Steidel, & Boksenberg (1988) measure C IV $`\lambda \lambda 1548`$, 1551 in the 1.486 and 1.524 systems with equivalent widths $`0.7`$ Å. However, Steidel & Sargent (1992) give a spectrum showing no detectable Mg II absorption.
Simple estimates suggested that the $`\lambda 2300`$ flux drop might be too gradual to be attributed to the converging Lyman lines of the $`z=1.486`$ system. Therefore, we computed a model spectrum involving an assumed power-law continuum and absorption by the hydrogen lines and bound-free continuum. The column density of H I was parameterized by the Lyman edge optical depth, $`\tau _H`$, and the lines were assumed to have a Gaussian profile with a Doppler parameter set to a typical value $`b=30\mathrm{km}\mathrm{s}^1.`$ (A larger line width would give an excessive equivalent width for L$`\alpha `$ for the required $`\tau _\mathrm{H}`$ . The observed decrement of the Lyman line equivalent widths actually suggests a narrower core line width, with some of the L$`\alpha `$ equivalent width being attributable to a broader component of modest column density. These details do not affect our conclusions.) Inclusion of 50 Lyman lines proved more than adequate to trace the convergence to the continuum optical depth. The model spectrum was convolved with a single component Gaussian instrumental line profile following the discussion of Impey et al. (1996), using a FWHM of 2.9 Å for G190H and 4.2 Å for G270H. If only the $`z=1.486`$ system contributes to the LLS, then the Lyman line convergence is too close to the Lyman limit ($`\lambda 912`$) to fit the observed feature. However, allocation of some H I column density to both the 1.486 and 1.524 systems gave a good fit. Two additional LLS appear to be associated with the $`z=1.938`$ and 1.174 systems. Figure 3 shows the resulting model spectrum, along with the observed flux, binned in intervals of $`7.3`$ Å. The model has $`\tau _\mathrm{H}=(1.0,0.5,0.8,0.6)`$ for z = (1.174, 1.486, 1.524, 1.938), respectively. (The value of $`\tau _\mathrm{H}`$ for z = 1.174 is uncertain because of the poorly determined amount of scattered light below 2000 Å.) The three distinguishable LLS show good agreement with the expected $`\nu ^3`$ behavior of the Lyman continuum optical depth above threshold, for an intrinsic continuum slope in the range $`\alpha 1.5`$ to -2.0. This is consistent with the slope -1.8 found by Zheng et al. (1997) for their composite QSO spectrum in the wavlength range $`\lambda 600`$ to $`\lambda 1050`$. A value $`\alpha =1.8`$ is assumed in the fit shown in Figure 3. The gradual descent of the observed flux toward the Lyman limit for the $`z=1.938`$ system is a puzzle, but it may involve the effects of unrelated absorption lines. An understanding of this is important, as it would play a role in the classification of PG 1222+228 as a candidate Lyman edge QSO. Observations at higher spectral resolution would help to clarify the situation.
Sargent, Steidel, and Boksenberg (1989) discuss the statistics of LLS in QSOs. For the redshift range in question, they give a mean incidence of LLS of $`N(z)1.5`$ per unit redshift. Their data show many QSOs with multiple LLS, although the number of LLS in PG 1222+228 may be somewhat higher than typical. However, the efforts to measure Lyman continuum polarization in QSOs to some extent targeted the candidate Lyman edge QSOs, and objects with LLS of moderate optical depth may have an enhanced probability to be included.
We conclude that an intrinsic power-law continuum, together with cosmologically intervening Lyman limit absorption, provides a straightforward explanation of the ultraviolet spectrum of PG 1222+228.
## 4 THE INTRINSIC POLARIZED CONTINUUM
The various LLS in PG 1222+228 substantially attentuate the observed continuum. What is the behavior of the polarized flux when corrected for the absorption? In view of the uncertainties in the measured polarization, a sufficient procedure is to estimate the polarized flux by multiplying the measured polarization in the chosen wavelength bins by the assumed intrinsic continuum flux. For this, we use the same wavelength bins described earlier, and the $`I_{\nu }^{}{}_{}{}^{PL}\nu ^{1.8}`$ power-law continuum used in our fit. The resulting rotated Stokes flux, $`Q_{}^{}=q^{}\times I_{\lambda }^{}{}_{}{}^{PL}`$, is shown in Figure 4. We see that the Stokes flux now rises strongly with decreasing wavelength in the region of the polarization rise. This resembles the result found for PG 1630+377 by Koratkar et al. (1995). A rising polarized flux is an important constraint on models for the origin of the polarization rise.
SWH showed that the wavelength dependence of the flux and polarization in PG 1222+228 and PG 1630+377 could be fit with an ad hoc model involving an accretion disk. The disk radiates as a black body, but the brightness temperature is depressed below the effective temperature for wavelengths below the Lyman limit, simulating a Lyman edge in the disk atmosphere. The polarization is assumed to rise abruptly at $`\lambda 912`$ by an arbitrary amount. Relativistic effects give a blueshifted, but still fairly abrupt polarization rise in the observed spectrum. The models are characterized by $`a_{},`$ the dimensionless angular momentum of the hole; the black hole mass; the accretion rate, $`\dot{M}_0\dot{M}/(1\mathrm{M}_{}\mathrm{yr}^1)`$; and the viewing angle, $`\mu _{obs}=cos(\theta _{obs})`$. SWH found that $`a_{}=0.5`$ gave approximately the observed wavelength for the polarization rise, for a relatively edge-on viewing angle. Their fits to both objects had a fairly low value of $`T_{max}`$, the maximum disk effective temperature, as required by the dropping flux in the Lyman continuum region. The correction for LLS absorption in PG 1222+228 hardens the far ultraviolet spectrum, and a higher value of $`T_{max}`$ is required to fit the energy distribution. Figure 5 shows the continuum flux and polarization for a model with $`a_{}=0.5,`$ $`M_9=8.8`$ and $`\dot{M}_0=86`$. (We have assumed H<sub>0</sub>= 70 $`\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$and q<sub>0</sub>= 0.5.) The model agrees reasonably well with the corrected flux in the Lyman continuum and with the longer wavelength measurements. Although the corrected flux was assumed to be a power law, a disk continuum would also likely fit the observed flux, given some freedom to adjust the LLS optical depths. This model has a step-function rise in polarization from negligle polarization at wavelengths longward of $`\lambda 912`$ to an ad hoc value of 2.1 times the Chandrasekhar (1960) value for a pure scattering atmosphere. The model predicts an observed polarization and polarized flux that rise at a wavelength substantially blueshifted from $`\lambda 912`$, but the rise is more gradual than observed. This underscores the need for improved polarization measurements of this object. If the relativistic transfer function gives too gradual a polarization rise, even for an instantaneous polarization rise in the rest frame of the orbiting gas, then accretion disk models for the polarization rise will face a serious problem.
The adopted model parameters give a bolometric luminosity $`L/L_E=0.36,`$ where $`L_E`$ is the Eddington limit. Such a high value of $`L/L_E`$ is barely consistent with a geometrically thin disk. A larger value of $`a_{}`$ would allow a larger mass for the required $`T_{max}`$, but then the polarization rise would occur at a wavelength shorter than observed (cf. SWH). Evidently, an accretion disk fit to the corrected continuum of PG 1222+228, in the manner of SWH, pushes the disk parameters to the limits. Conceivably, the thickening of the disk corresponding to the large value of $`L/L_E`$ may be related to the origin of the polarization rise.
## 5 DISCUSSION
The Lyman continuum polarization rises are among the more puzzling recent observational discoveries concerning QSOs. The wavelength dependence, rising rather abruptly at nearly the same rest wavelength in the several known cases, suggests a connection with the bound levels of atoms. The proximity of the feature to $`\lambda 912`$ further suggests an association with the Lyman edge of hydrogen. The ad hoc model of SWH supports an association with the Lyman edge and raises the possiblity of confirming the presence of a relativistic disk and constraining its parameters. However, attempts to fit the feature with a physical model have encountered difficulties. This is an important problem for QSO theory.
Are the reported polarization rises real? The coincidence of the polarization rise in PG 1222+228 with a sharp drop in flux might raise the question of background problems with the FOS spectropolarimeter. Impey et al. (1995) argue that the polarization is unlikely to be less than 2.7% around 2000 Å under any reasonable assumption for the FOS background. However, the degree of polarization in the far ultraviolet is uncertain by at least a factor two because of systematic errors involving background and scattered light in the FOS. The observed polarization rises in several QSOs occur at different observed wavelengths but similar rest wavelengths. We are not aware of any polarization rises of this nature in FOS spectropolarimetry of BL Lac objects or stars. There is an urgent need for a renewed capability for ultraviolet spectropolarimetry from space to confirm and extend the measurements.
Lyman continuum polarization rises have heretofore been associated with the candidate Lyman edge QSOs (see discussion by Koratkar et al. 1998). Our results suggest that PG 1222+228 may not be a true member of this class. Measurements of the Lyman continuum polarization in additional QSOs are needed to clarify the frequency of occurence of the phenomenon and to look for correlations with features in the continuum flux, in the line intensities and polarization, and other properties. Observations to shorter rest wavelengths are needed to determine whether the polarization falls or continues to rise. Measurements of the time dependence of the polarization rises would be most interesting. The emitting radius of an accretion disk would be light weeks.
Lee and Blandford (1995) considered a model for the far ultraviolet polarization rise of QSOs that did not involve the Lyman edge. Noting that a number of resonance lines of heavy elements fall in the rest wavelength range where the polarization rises, they suggested that resonance scattering of the QSO continuum might produce the observed polarization. Such a model could produce a rising polarized flux, since the polarization could be essentially zero at wavelengths without scattering contributions. As noted above, the polarization rise in PG 1222+228 may be too steep for models involving an accretion disk. In this case, alternative models such as resonance scattering may hold promise. We note that the polarized flux spectrum of PG 1630+377 (Koratkar et al. 1995) shows a strong rise at the wavelength of the N V emission line.
The claim by Richards et al. (1999) that luminous QSOs have many intrinsic, narrow, high velocity C IV absorption systems has important implications. This would complicate the use of such systems to probe the evolution of galaxies and the intergalactic medium. The mechanism for producing the absorbing clouds would add another challenge to the subject of outflows from QSOs. Surveys of QSOs at different luminosities, to a uniform standard of signal-to-noise ratio, would allow confirmation of the claimed higher incidence of absorption in more luminous QSOs. Tests of the intrinsic nature of narrow absorptions have been summarized by Barlow, Hamann, & Sargent (1997). These include time variability of the depth and profile of the lines and evidence for saturated but nonblack line profiles. Detection of changes in the velocities of the lines would be most revealing. Chemical abundances may be an indicator, given evidence for high abundances in the broad absorption and emission-line regions of QSOs (e.g., Hamann & Ferland 1993). If C IV systems in high and low luminosity QSOs have similar, mostly subsolar abundances, this might argue for their intervening nature.
###### Acknowledgements.
The author is grateful to C. Impey and C. Petri for providing advice and the reduced HST observations and to C. Sneden for the use of a computer subroutine. The work has benefitted from discussions and communications with R. Antonucci, O. Blaes, R. Blandford, R. Ganguly, M. Malkan, R. Mushotzky, G. Richards, R. Weymann, and B. Wills. This material is based in part upon work supported by the Space Telescope Science Institute under Grant No. GO-07359.02. This work was carried out in part at the Institute for Theoretical Physics, University of California, Santa Barbara, supported by NSF grant PH94-07194.
Captions for Figures |
warning/0002/math0002068.html | ar5iv | text | # Metastability of Breather Modes of Time Dependent Potentials
## 1 Introduction and Overview
We are interested in the initial value problem for the linear Schrödinger equation in one space dimension
$$i_tf=\left(\frac{1}{2}_x^2+V(x,t)\right)f(t)f$$
(1.1)
Here, $`f`$ is a complex-valued function of $`x`$ and $`t`$. We assume that $`V(x,t)`$ is a smooth real-valued potential energy function which is sufficiently localized in space (for example, of Schwartz class). In our specific applications, $`V(x,t)`$ will be taken to be a periodic function of $`t`$. However, the techniques we use can be adapted for to more general time-dependence . Note that (1.1) is a nonautonomous Hamiltonian system:
$$i_tf=\frac{\delta h}{\delta f^{}}[f,f^{},t],$$
where $`h`$ denotes the Hamiltonian energy:
$$h[f,f^{},t]=\left(\frac{1}{2}|_xf|^2+V(x,t)|f|^2\right)dx.$$
If $`V(x,t)`$ is not independent of $`t`$, $`h`$ is not a conserved integral of the flow. On the other hand, since the potential $`V`$ is real-valued, the flow defined by (1.1) is always unitary in $`L^2()`$, i.e.
$$|f(x,t)|^2𝑑x=|f(x,0)|^2𝑑x,t.$$
(1.2)
In applications, it is often natural to decompose $`V(x,t)`$ as:
$`V(x,t)=V_0(x,t)+W(x,t),`$
where $`V_0(x,t)`$ denotes an unperturbed potential, and $`W(x,t)V(x,t)V_0(x,t)`$, denotes a small perturbation. Thus,
$`(t)_0(t)+𝒲(t),`$ (1.3)
and (1.1) can be rewritten as:
$$i_tf=\left(_0(t)+𝒲(t)\right)f$$
(1.4)
Here, we have denoted the multiplication operator $`fW(x,t)f`$ by $`𝒲(t)`$. The choice of $`V_0(x,t)`$ is often dictated by some a priori knowledge of the solutions of the unperturbed system
$`i_tf=_0(t)f.`$ (1.5)
A problem of importance is then to contrast the detailed dynamics of solutions to (1.4) with those of the unperturbed system (1.5). In particular, if (1.5) has bound state solutions (breather modes, or solutions having finite energy and not decaying as $`|t|\mathrm{}`$) do they persist in the perturbed dynamical system (1.4)?
The simplest variant of this problem is the case where the unperturbed part is stationary, i.e. $`V_0(x,t)=V_0(x)`$. Suppose the operator $`_0`$ has an $`L^2`$ eigenfunction. The unitary evolution of the spatially localized eigenfunction is time-periodic and represents a bound state solution of the unperturbed Schrödinger equation (1.5). The perturbed model (in this and in the more general case when $`_0`$ has multiple discrete eigenvalues) is related to the problem of ionization of an atom by a time-dependent electromagnetic field and the problem of describing the effects of weak inhomogeneities on the propagation of continuous waves in optical fibers . Using a time-dependent method developed in the context of (i) quantum resonances and the perturbation theory of embedded eigenvalues in the continuous spectrum and (ii) resonances and radiation damping of bound nonlinear wave equations , Soffer and Weinstein studied the metastability of such states . Generalizations of this theory for handling multifrequency perturbations and the interference of multiple bound states in the unperturbed problem have been explored by Kirr and Weinstein. Based on the observation that the mechanism for instability of the bound state is coupling of the bound state to the continuous spectral modes, the analysis was carried out at the level of the coupled equations for the bound state and dispersive components of the solution. Under general hypotheses on unperturbed Hamiltonian (local energy decay estimates on the unitary propagator $`e^{it_0}`$) this equivalent dynamical system was studied and it was shown that a bound state is generically unstable but long-lived. The lifetime is given by a formula analogous to the Fermi golden rule .
In this paper we consider the case where the unperturbed Hamiltonian is genuinely time-dependent. A physical application of the theory we develop, in the context of frequency detuning in periodically modulated optical waveguides , will be presented in §5. Let $`_0(t)=\frac{1}{2}_x^2+V_0(x,t)`$, where $`V_0(x,t)`$ is smooth, periodic in $`t`$ with the same period for each $`x`$ and of sufficiently rapid decay for large $`x`$ for each $`t`$. The particular choices of $`V_0(x,t)`$ we consider in this paper belong to a large family of very special, so-called separable, time-dependent potentials, $`V_0(x,t),x`$, studied by Miller and Akhmediev . The separable potential $`V_0(x,t)`$ can be chosen to be time-periodic, in which case the unperturbed problem supports exact bound states (breather modes) and the initial value problem for (1.5) can be solved exactly. That is, a complete set of eigenmodes and generalized eigenmodes can be explicitly displayed with respect to which the dynamics of (1.5) is diagonal. This class of potentials is intimately connected with the soliton theory of completely integrable multicomponent cubic nonlinear Schrödinger equations .
The existence of such exact breather modes in the unperturbed time-periodic problem is quite remarkable and we believe that this is a highly non-generic phenomenon<sup>1</sup><sup>1</sup>1The scarcity of breather solutions of nonlinear wave equations defined on a spatial continuum of infinite extent has been extensively explored in the setting of perturbations of the completely integrable sine-Gordon equation; see, for example, . The connection with linear nonautonomous problems can be made by viewing a breather solution of a nonlinear dynamical problem as a bound state of a linear problem with a given (self-consistent) potential.. Indeed, from a general dynamical systems perspective, (1.5) with such a choice of $`V_0(x,t)`$, may be viewed as a parametrically forced wave equation (here we are actually considering the time-periodic function $`V_0(x,t)`$ itself to be the sum of a time-independent part and a time-periodic modulation). One therefore expects that the presence of resonances will perturb the Floquet multipliers (corresponding to bound states) off the unit circle as in the elementary example of Mathieu’s equation . The persistence of breather solutions under the time-periodic perturbation would imply the non-departure of a Floquet multiplier from the unit circle to all orders in the size of the perturbation. The fact that infinitely many such conditions hold for these special separable potentials is no doubt linked to the infinite sequence of symmetries and time invariants enjoyed by the completely integrable nonlinear flow that underpins the construction of the separable potentials (see Appendix A for more details). Of course, this is only a heuristic picture. In fact, the perturbation theory of the Floquet multipliers is complicated by the fact that they are embedded in the continuous spectrum which covers the unit circle. However, spectral deformation methods have been developed for some classes of models that could well be adapted here. Relevant technical details can be found in <sup>2</sup><sup>2</sup>2 For nonlinear wave equations defined on an infinite lattice (e.g. discrete sine-Gordon, discrete $`\varphi ^4`$), breather solutions can be constructed for sufficiently large lattice spacing; see for example . The radiative decay of such discrete breathers, for sufficiently small lattice spacing, is expected to be governed by a mechanism of the kind studied in this paper; see also . Related to this are results for the dynamics of kinks of discrete nonlinear wave equations, in which the techniques of this paper have been used to study the “pinning” of discrete kinks on a lattice site. This pinning is marked by the slow radiative decay of spatially localized and time-periodic or quasiperiodic oscillations about a static kink . .
We want to make our motivation for pursuing deformations of these admittedly rather special periodic potentials very clear. First of all, the problem is relevant to the analysis of optical waveguides. In the paraxial approximation, the slowly varying envelope of a highly oscillatory electric field in a dielectric medium with inhomogeneous dielectric properties (index of refraction) satisfies an equation of form (1.1). Here, $`t`$ denotes the longitudinal variable, the direction of propagation, and $`x`$ is the transverse spatial variable. For inhomogeneous index profiles corresponding to exactly solvable potentials, light of a particular frequency propagates as a non-attenuating bound state mode in these wave guides. However, if the light frequency deviates from the “integrable frequency” the propagating wave will be governed by the perturbed equation (1.4). Thus the question of whether such modes persist and if not what their lifetime is for the perturbed dynamics naturally arises. We will give more details about this problem in §5.
But it is also true that the study of perturbed separable periodic potentials is important in general terms. Given an arbitrary time-periodic potential in the Schrödinger equation, one wants to study the corresponding dynamics using perturbation theory. In doing so, the first question one must address is that of finding a “nearby” problem that can be solved exactly. We simply take the point of view that many periodic potentials will be closer to a separable periodic potential (in a sense that can be made precise) than to any time-independent potential.
In any case, with the explicit spectral theory associated with $`V_0(x,t)`$ in hand, our goal is to carry out a detailed analytical study of the coupled mode dynamics induced by a time-dependent perturbation $`W(x,t)`$. We establish the generic metastabilty of the exact bound states associated with separable periodic potentials $`V_0(x,t)`$ and obtain a detailed picture of the dynamics.
The paper is structured as follows. In §2 we first review the construction of time-dependent exactly solvable potentials from a set of discrete data, and then show how the initial value problem for such separable potentials can be solved explicitly. We then describe how properties of the separable potentials depend on the choice of the discrete data generating them. Next, we derive by projection onto an orthonormal basis the general coupled mode equations which arise when a separable potential is perturbed by some arbitrary correction $`W(x,t)`$. This section will then conclude with a detailed derivation of the two-soliton time-periodic potential and its associated explicit complete set of bound states and generalized eigenfunctions. More details about the separable potentials described in §2 are given in Appendix A.
In §3 Floquet theory is then used to map the coupled mode equations to a system associated with a time-dependent perturbation of an autonomous system, a situation analyzed in detail in and . In §4, we then describe the dynamics of solutions of the coupled mode equations for even time-periodic perturbations $`W(x,t)`$ of a separable two-soliton even time-periodic potential $`V_0(x,t)`$ ($`V_0(x,t)`$ and $`W(x,t)`$ both share even parity in $`x`$ and have the same temporal period). In particular, we study the initial value problem when the initial condition is a pure bound state of the unperturbed problem. First, we study the small time behavior of the coupled mode equations (without requiring $`W(x,t)`$ to be small) and deduce that the bound state amplitude behaves as $`1Ct^2`$ for some constant $`C`$ and interpret this result in the context of the theory of ideal measurements in quantum mechanics (the “watched pot” effect). We then assume the perturbation $`W(x,t)`$ to be small, of size $`ϵ`$, and seek the behavior of the bound state amplitude over intermediate times of order $`ϵ^1`$ and $`ϵ^2`$ using the classical method of multiple scales.
We show the existence of a perturbation-induced frequency shift of the breather mode evident on timescales of order $`ϵ^1`$ and exponential decay of the bound state mode amplitude on timescales of order $`ϵ^2`$. The condition for the decay constant to be nonzero is a direct analog of the “Fermi golden rule”.
Then, using the transformation to an autonomous system found in §3, we show how the rigorous theory developed for multifrequency perturbations of autonomous systems by Kirr and Weinstein can be applied in some cases to justify the multiple scales calculation, and to provide more detailed information about the infinite time behavior of the solution. This analysis completes the portrait of the dynamics, showing that the exponential decay is ultimately washed out in a sea of dispersive waves, at which point the decay becomes algebraic in time.
Having described the theory, in §5 we consider an application of the analysis to a problem of frequency detuning in planar optical waveguides. Finally, in §6 the prediction of an exponential decay constant $`\mathrm{\Gamma }`$ for the bound state mode amplitude found in section 4 is compared to numerical simulations of the perturbed time-dependent Schrödinger equation.
A detailed description of the theory of separable potentials, at once summarizing for completeness and also further developing the results of , can be found in Appendix A. In Appendix B the reader will find the proofs of the decay estimates that we will use in §4 in order to apply the results of .
Regarding notation: We will use the inner product
$$f(),g()=_{\mathrm{}}^{\mathrm{}}f(x)^{}g(x)𝑑x$$
(1.6)
on $`L^2()`$. Occasionally, the angled brackets will denote the inner product in more general Hilbert spaces. Linear operators will be denoted with calligraphic letters, vectors with arrows, and matrices with boldface letters. We will often use the function defined by:
$$x(1+x^2)^{1/2}.$$
(1.7)
Complex conjugation will be denoted with stars, and time averages will be denoted with bars.
## 2 Exactly solvable time-dependent potentials
In this section we recall for our purposes a class of time-dependent potentials $`V_0(x,t)`$ related to $`M`$-soliton solutions of certain completely integrable nonlinear flows. Because of the intimate connection of these potentials to integrable systems, it is possible to explicitly derive the spectral representation associated with such potentials . This section is divided into five parts. First, the direct construction of separable potentials from a set of discrete data $`𝒟`$ is outlined. Then, we show how the same discrete data $`𝒟`$ gives rise to formulas for a complete set of modes for the time-dependent Schrödinger equation corresponding to the separable potential $`V_0(x,t)`$ and how this basis is easily used to express the general solution of the initial value problem. We then give a qualitative description of the kinds of functions $`V_0(x,t)`$ one can obtain from this procedure. As we ultimately want to consider perturbations of $`V_0(x,t)`$, we next show how to use the basis of solutions to the unperturbed problem to derive the coupled mode equations which trivialize the unperturbed dynamics and lay bare the perturbative effects. Finally, we specialize to the case of an even periodic potential corresponding to a two-soliton solution of the cubic nonlinear Schrödinger equation. As one might anticipate, the evenness (in $`x`$) of the potential leads to some simplifications in the spectral representation.
### 2.1 Separable time-dependent potentials.
Let us present the construction of the family of time-dependent potentials that we will consider in this paper, and describe their properties with respect to the linear Schrödinger equation. More details can be found in Appendix A.
Each potential we shall consider will be specified by a certain set of discrete data. Let $`N`$ and $`M`$ be independent natural numbers. A set of discrete data $`𝒟`$ consists of $`M`$ distinct complex numbers $`\lambda _1,\mathrm{},\lambda _M`$ in the upper half-plane, and $`M`$ vectors $`\stackrel{}{g}^{(1)},\mathrm{},\stackrel{}{g}^{(M)}`$ in $`^N`$.
The discrete data $`𝒟`$ is used to build a potential function $`V_0(x,t)`$ in the following way. Introduce the scalar expression
$$a(x,t,\lambda )=\left(\lambda ^M+\underset{p=0}{\overset{M1}{}}\lambda ^pa^{(p)}(x,t)\right)e^{2i(\lambda x+\lambda ^2t)},$$
(2.1)
and the $`N`$-component vector expression
$$\stackrel{}{b}(x,t,\lambda )=\underset{p=0}{\overset{M1}{}}\lambda ^p\stackrel{}{b}^{(p)}(x,t).$$
(2.2)
In these expressions, the coefficients $`a^{(p)}(x,t)`$ and $`\stackrel{}{b}^{(p)}(x,t)`$ are undetermined functions of $`x`$ and $`t`$. They will now be determined with the use of the discrete data $`𝒟`$. For $`k=1,\mathrm{},M`$, we insist that $`a(x,t,\lambda )`$ and $`\stackrel{}{b}(x,t,\lambda )`$ satisfy the relations:
$$\begin{array}{ccc}\hfill a(x,t,\lambda _k)& =& \stackrel{}{g}^{(k)}\stackrel{}{b}(x,t,\lambda _k),\hfill \\ \hfill \stackrel{}{b}(x,t,\lambda _k^{})& =& a(x,t,\lambda _k^{})\stackrel{}{g}^{(k)}.\hfill \end{array}$$
(2.3)
These equations amount to a square linear inhomogeneous system of algebraic equations for the coefficient functions $`a^{(p)}(x,t)`$ and the components of $`\stackrel{}{b}^{(p)}(x,t)`$. We will soon illustrate this procedure with a concrete example. From the solution of this linear system, the potential function connected with the discrete data $`𝒟`$ is given in terms of the components of $`\stackrel{}{b}^{(M1)}(x,t)`$ by
$$V_0(x,t)4\underset{n=1}{\overset{N}{}}\left|b_n^{(M1)}(x,t)\right|^2.$$
(2.4)
This function $`V_0(x,t)`$ is a genuinely time-dependent potential well. Furthermore, it can be shown that $`V_0(x,t)`$ is in the Schwartz space as a function of $`x`$, and its $`L^1`$ norm is constant in $`t`$. The latter follows from the fact that $`V_0(x,t)`$ can be viewed as the self-consistent nonlinear potential for an $`N`$-component cubic nonlinear Schrödinger equation, which conserves the sum of the $`L^2`$ norms of the $`N`$ field components, which are proportional to the $`b_n^{(M1)}(x,t)`$ for $`n=1,\mathrm{},N`$.
### 2.2 Solution of the linear Schrödinger equation with a separable potential.
Along with the potential function $`V_0(x,t)`$, this construction starting from the discrete data $`𝒟`$ also provides all of the solutions of the corresponding linear Schrödinger equation . These are built from the function $`a(x,t,\lambda )`$ as follows. For all real $`\lambda `$, set
$$\mathrm{\Psi }_\mathrm{d}(x,t,\lambda )\left(\pi \underset{k=1}{\overset{M}{}}|\lambda \lambda _k|^2\right)^{1/2}a(x,t,\lambda ),$$
(2.5)
and let the functions $`\mathrm{\Psi }_{\mathrm{b},1}(x,t),\mathrm{},\mathrm{\Psi }_{\mathrm{b},M}(x,t)`$ be the result of applying the Gram-Schmidt procedure (in $`L^2()`$ with respect to $`x`$) to the functions $`a(x,t,\lambda _1^{}),\mathrm{},a(x,t,\lambda _M^{})`$ at any fixed value of $`t`$. Then we have :
1. Each function $`\mathrm{\Psi }_\mathrm{d}(x,t,\lambda )`$ for $`\lambda `$ and each function $`\mathrm{\Psi }_{\mathrm{b},k}(x,t)`$ is a solution of the linear Schrödinger equation with potential $`V_0(x,t)`$. The fact that the $`L^2`$ inner product is an invariant of the evolution shows that the functions $`\mathrm{\Psi }_{\mathrm{b},k}(x,t)`$ do not depend on the choice of the time $`t`$ at which they are obtained from the Gram-Schmidt process.
2. For any fixed $`t`$, these functions form an orthonormal basis of $`L^2()`$.
These facts show us how to solve the initial value problem for the linear Schrödinger equation for the potential $`V_0(x,t)`$. Namely, to find the solution of
$$i_tf+\frac{1}{2}_x^2fV_0(x,t)f=0,f(x,0)=f_0(x)L^2(),$$
(2.6)
one simply projects the initial data onto the basis at $`t=0`$ by defining:
$$\widehat{f}(\lambda )\mathrm{\Psi }_\mathrm{d}(,0,\lambda ),f_0(),\widehat{f}_k\mathrm{\Psi }_{\mathrm{b},k}(,0),f_0(),$$
(2.7)
and then recovers the solution as a superposition of modes:
$$f(x,t)=\underset{k=1}{\overset{M}{}}\widehat{f}_k\mathrm{\Psi }_{\mathrm{b},k}(x,t)+_{\mathrm{}}^{\mathrm{}}\widehat{f}(\lambda )\mathrm{\Psi }_\mathrm{d}(x,t,\lambda )𝑑\lambda .$$
(2.8)
If the potential $`V_0(x,t)`$ is slightly perturbed, it may still be convenient to expand in this basis, but then the coefficients $`\widehat{f}(\lambda )`$ and $`\widehat{f}_1,\mathrm{},\widehat{f}_M`$ will become time-dependent.
### 2.3 Qualitative description of separable potentials.
Let us describe the types of potential functions $`V_0(x,t)`$ that can be obtained by this procedure. In the generic case when the real parts of the parameters $`\lambda _1,\mathrm{},\lambda _M`$ are all distinct, these potentials have the form of a collision among $`M`$ moving potential wells. That is, as $`t\pm \mathrm{}`$,
$$V_0(x,t)\underset{k=1}{\overset{M}{}}V_0^{(k)\pm }(x,t),$$
(2.9)
where the individual wells have the form
$$V_0^{(k)\pm }(x,t)=4\rho _k^2\mathrm{sech}^2(2\rho _k(x+2\sigma _kt)\delta _k^\pm ),$$
(2.10)
where $`\delta _k^\pm `$ are constants that depend on the vectors $`\stackrel{}{g}^{(1)},\mathrm{},\stackrel{}{g}^{(M)}`$, and where $`\lambda _k=\sigma _k+i\rho _k`$. Considered in isolation, each well carries a single bound state. When the wells collide for finite $`t`$, the bound states can become mixed, and a state $`f(x,t)`$ that is bound in a single well as $`t\mathrm{}`$ will have a component in each well as $`t+\mathrm{}`$. The associated scattering matrix can be computed exactly .
If some of the parameters $`\lambda _k`$ share the same real part $`\sigma `$, then the asymptotics of the potential $`V_0(x,t)`$ in the frame moving with velocity $`2\sigma `$ will no longer be stationary, but will be generally quasiperiodic. In particular, if all of the parameters $`\lambda _k`$ are purely imaginary, then the potential $`V_0(x,t)`$ will generally be a quasiperiodic function of the time $`t`$. This is clear because taking $`\lambda _k=i\rho _k`$ with $`\rho _k`$ real and positive ensures that the only time dependence that enters into the computation of $`V_0(x,t)`$ is via the exponentials $`\mathrm{exp}(\pm 2i\rho _k^2t)`$. Such a potential is automatically quasiperiodic. We can further ensure that the potential function is strictly periodic by making the frequencies commensurate. This will be true<sup>3</sup><sup>3</sup>3For $`M>2`$. The potential is always periodic if $`M=2`$ and is stationary if $`M=1`$. if the parameters $`\rho _k`$ have the form
$$\rho _k=\sqrt{n_k\frac{\mathrm{\Omega }_0}{2}+\mathrm{\Delta }},$$
(2.11)
where $`\mathrm{\Omega }_0`$ is some fundamental frequency and $`n_k`$ are distinct integers. This choice ensures that the frequencies
$$\omega _{jk}2\rho _j^22\rho _k^2=(n_jn_k)\mathrm{\Omega }_0$$
(2.12)
are all integer multiples of $`\mathrm{\Omega }_0`$. Only the frequency differences $`\omega _{jk}`$ are important because the potential is given as a sum of absolute values (2.4). In fact, it can be seen from the form of the linear system (2.3) that
$$b_n^{(p)}(x,t)=e^{2i(\rho _1^2+\mathrm{}+\rho _M^2)t}G_{n,p}(\{e^{i\omega _{jk}t}\},x),$$
(2.13)
where $`G_{n,p}`$ is, for each fixed $`x`$, a rational function of the exponentials $`\mathrm{exp}(i\omega _{jk}t)`$. The sufficiency of the relations (2.11) to guarantee time periodicity of $`V(x,t)`$ with fundamental frequency $`\mathrm{\Omega }_0`$ is then clear from (2.4).
### 2.4 Perturbed separable potentials and coupled mode equations.
As we have already suggested, the explicit basis of exact solutions derived in the previous subsection forms a natural coordinate system in which to study perturbed problems. Let $`W(x,t)`$ be a correction to the potential energy, so that the equation becomes
$$if_t=\left(\frac{1}{2}_x^2+V_0(x,t)\right)f+W(x,t)f=_0(t)f+𝒲(t)f.$$
(2.14)
Here, $`V_0(x,t)`$ is a separable time-dependent potential built from the discrete data $`𝒟=\{\lambda _1,\mathrm{},\lambda _M,\stackrel{}{g}^{(1)},\mathrm{},\stackrel{}{g}^{(M)}\}`$. So we use completeness to express $`f(x,t)`$ for each fixed $`t`$ in terms of the basis of solutions of the unperturbed problem:
$$f(x,t)=\underset{k=1}{}B_{\mathrm{b},k}(t)\mathrm{\Psi }_{\mathrm{b},k}(x,t)+_{\mathrm{}}^{\mathrm{}}B_\mathrm{d}(t,\lambda )\mathrm{\Psi }_\mathrm{d}(x,t,\lambda )𝑑\lambda .$$
(2.15)
In the absence the the perturbation $`W(x,t)`$, the mode amplitudes $`B_{\mathrm{b},k}`$ and $`B_\mathrm{d}(t,\lambda ),\lambda `$ are governed by the equations $`_tB_{\mathrm{b},k}=0,_tB_\mathrm{d}(t,\lambda )=0`$. In the presence of a perturbation $`W(x,t)`$ coupled mode equations can be derived by projecting (2.14) onto the basis elements $`\mathrm{\Psi }_{\mathrm{b},k}(x,t)`$ and $`\mathrm{\Psi }_\mathrm{d}(x,t,\lambda )`$. This yields the system of coupled mode equations:
$$\begin{array}{ccc}\hfill i_t\stackrel{}{B}_\mathrm{b}(t)& =& 𝐌(t)\stackrel{}{B}_\mathrm{b}(t)+_{\mathrm{}}^{\mathrm{}}B_\mathrm{d}(t,\lambda )\stackrel{}{N}(t,\lambda )𝑑\lambda ,\hfill \\ & & \\ \hfill i_tB_\mathrm{d}(t,\eta )& =& \stackrel{}{N}(t,\eta )^{}\stackrel{}{B}_\mathrm{b}(t)+_{\mathrm{}}^{\mathrm{}}K(t,\eta ,\lambda )B_\mathrm{d}(t,\lambda )𝑑\lambda ,\hfill \end{array}$$
(2.16)
for the coefficients of $`f(x,t)`$, where $`\stackrel{}{B}_\mathrm{b}(t)`$ is the vector of bound state amplitudes $`B_{\mathrm{b},k}(t)`$, and where the matrix elements of the perturbation $`W(x,t)`$ are explicitly given by
$$\begin{array}{ccc}\hfill M_{kj}(t)& =& \mathrm{\Psi }_{\mathrm{b},k}(,t),W(,t)\mathrm{\Psi }_{\mathrm{b},j}(,t),\hfill \\ & & \\ \hfill N_k(t,\lambda )& =& \mathrm{\Psi }_{\mathrm{b},k}(,t),W(,t)\mathrm{\Psi }_\mathrm{d}(,t,\lambda ),\hfill \\ & & \\ \hfill K(t,\eta ,\lambda )& =& \mathrm{\Psi }_\mathrm{d}(,t,\eta ),W(,t)\mathrm{\Psi }_\mathrm{d}(,t,\lambda ),\hfill \end{array}$$
(2.17)
where $`N_k(t,\lambda )`$ are the components of the vector $`\stackrel{}{N}(t,\lambda )`$ and $`M_{kj}(t)`$ are the elements of the matrix $`𝐌(t)`$. In particular, it follows that the matrix $`𝐌(t)`$ is Hermitian and the scalar kernel $`K(t,\eta ,\lambda )`$ is Hermitian symmetric because $`W(x,t)`$ is real. With the unperturbed problem exactly diagonalized in this way, this system is a useful starting point for perturbation theory.
### 2.5 Even two-soliton periodic potentials.
In this subsection, we illustrate the procedures described above in some detail with an example that is important in applications and that will guide the subsequent discussion. We consider the case $`N=1`$ and $`M=2`$, and accordingly introduce the expressions
$$a(x,t,\lambda )=(\lambda ^2+a^{(1)}(x,t)\lambda +a^{(0)}(x,t))e^{2i(\lambda x+\lambda ^2t)},b(x,t,\lambda )=b^{(1)}(x,t)\lambda +b^{(0)}(x,t).$$
(2.18)
Because $`N=1`$ these are both scalar expressions, and we have at the moment four complex valued unknown functions, $`a^{(0)}(x,t)`$, $`a^{(1)}(x,t)`$, $`b^{(0)}(x,t)`$ and $`b^{(1)}(x,t)`$. To find these, we introduce the discrete data $`\lambda _1`$, $`\lambda _2`$, $`g^{(1)}`$, and $`g^{(2)}`$ (again, here the $`g^{(k)}`$ are complex scalars because $`N=1`$). The linear equations (2.3) then become
$$\begin{array}{ccc}\hfill (\lambda _1^2+a^{(1)}(x,t)\lambda _1+a^{(0)}(x,t))e^{2i(\lambda _1x+\lambda _1^2t)}& =& g^{(1)}(b^{(1)}(x,t)\lambda _1+b^{(0)}(x,t))\hfill \\ & & \\ \hfill (\lambda _2^2+a^{(1)}(x,t)\lambda _2+a^{(0)}(x,t))e^{2i(\lambda _2x+\lambda _2^2t)}& =& g^{(2)}(b^{(1)}(x,t)\lambda _2+b^{(0)}(x,t))\hfill \end{array}$$
(2.19)
and
$$\begin{array}{ccc}\hfill b^{(1)}(x,t)\lambda _1^{}+b^{(0)}(x,t)& =& g^{(1)}(\lambda _1^2+a^{(1)}(x,t)\lambda _1^{}+a^{(0)}(x,t))e^{2i(\lambda _1^{}x+\lambda _1^2t)}\hfill \\ & & \\ \hfill b^{(1)}(x,t)\lambda _2^{}+b^{(0)}(x,t)& =& g^{(2)}(\lambda _2^2+a^{(1)}(x,t)\lambda _2^{}+a^{(0)}(x,t))e^{2i(\lambda _2^{}x+\lambda _2^2t)}.\hfill \end{array}$$
(2.20)
Given the discrete data $`𝒟`$, one can solve these equations for $`a^{(0)}(x,t)`$, $`a^{(1)}(x,t)`$, $`b^{(0)}(x,t)`$ and $`b^{(1)}(x,t)`$, say by Cramer’s rule, and thus obtain explicit expressions in terms of exponential functions.
Specializing to the case of $`\lambda _1=i\rho _1`$, $`\lambda _2=i\rho _2`$ (we assume without loss of generality that $`\rho _2>\rho _1`$), we obtain a time-periodic potential function, since the parameters $`\lambda _k`$ are pure imaginary and then the commensurability condition is automatically satisfied for $`M=2`$. Furthermore choosing $`g^{(1)}=e^{i\theta _1}`$ and $`g^{(2)}=e^{i\theta _2}`$ ensures that the potential function is even in $`x`$. Indeed, we then find that with $`s=\rho _2+\rho _1`$ and $`d=\rho _2\rho _1`$,
$$b^{(1)}(x,t)=2sd\frac{\rho _1\mathrm{cosh}(2\rho _2x)e^{2i\rho _1^2t+i\theta _1}\rho _2\mathrm{cosh}(2\rho _1x)e^{2i\rho _2^2t+i\theta _2}}{d^2\mathrm{cosh}(2sx)+s^2\mathrm{cosh}(2dx)4\rho _1\rho _2\mathrm{cos}(2sdt+\theta _2\theta _1)}.$$
(2.21)
The potential function is then given by
$$V_0(x,t)=4|b^{(1)}(x,t)|^2,$$
(2.22)
which is easily seen to be periodic in $`t`$ with period $`L=\pi /(sd)`$, and an even function of $`x`$. The shapes of these time-periodic potential wells are shown in Figures 1 and 2 for $`\theta _1=\theta _2=0`$ and two different choices of the parameters $`\rho _1`$ and $`\rho _2`$.
From the solution of the same linear system, we also find
$$a^{(0)}(x,t)=\rho _1\rho _2\frac{(\rho _1+\rho _2)^2S_1S_2+\rho _1^2e^{i\omega ti(\theta _2\theta _1)}+\rho _2^2e^{i\omega t+i(\theta _2\theta _1)}2\rho _1\rho _2C}{2\rho _1\rho _2\mathrm{cos}(\omega t+(\theta _2\theta _1))(\rho _1^2+\rho _2^2)C+(\rho _1+\rho _2)^2S_1S_2},$$
(2.23)
and
$$a^{(1)}(x,t)=i\frac{(\rho _1^2\rho _2^2)\rho _1C_2S_1+(\rho _2^2\rho _1^2)\rho _2C_1S_2}{2\rho _1\rho _2\mathrm{cos}(\omega t+(\theta _2\theta _1))(\rho _1^2+\rho _2^2)C+(\rho _1+\rho _2)^2S_1S_2},$$
(2.24)
where $`S_k\mathrm{sinh}(2\rho _kx)`$, $`C_k\mathrm{cosh}(2\rho _kx)`$, and $`C\mathrm{cosh}(2(\rho _1+\rho _2)x)`$, and where the frequency is $`\omega =2\pi /L=2sd`$. Note that $`a^{(0)}(x,t)`$ is an even function of $`x`$, while $`a^{(1)}(x,t)`$ is odd. We may then write the mode function $`a(x,t,\lambda )`$ in the form
$$\begin{array}{ccc}\hfill a(x,t,\lambda )& =& \left((\lambda ^2+a^{(0)}(x,t))\mathrm{cosh}(2i\lambda x)+\lambda a^{(1)}(x,t)\mathrm{sinh}(2i\lambda x)\right)e^{2i\lambda ^2t}+\hfill \\ & & \\ & & \left((\lambda ^2+a^{(0)}(x,t))\mathrm{sinh}(2i\lambda x)+\lambda a^{(1)}(x,t)\mathrm{cosh}(2i\lambda x)\right)e^{2i\lambda ^2t},\hfill \end{array}$$
(2.25)
in which the first term is even in $`x`$ and the second term is odd in $`x`$. Also, it is clear that $`a(x,t,\lambda )=a(x,t,\lambda )`$.
A particularly convenient orthonormal basis of the two-dimensional space of bound states is given by the formulas
$$\begin{array}{ccc}\hfill \mathrm{\Psi }_\mathrm{b}^{(\mathrm{e})}(x,t)& =& \frac{1}{\sqrt{4(\rho _1+\rho _2)}}\left[\frac{2}{\rho _1\rho _2}a(x,t,i\rho _1)+\frac{2}{\rho _2\rho _1}a(x,t,i\rho _2)\right],\hfill \\ & & \\ \hfill \mathrm{\Psi }_\mathrm{b}^{(\mathrm{o})}(x,t)& =& \frac{1}{\sqrt{4(\rho _1+\rho _2)}}\left[\sqrt{\frac{\rho _2}{\rho _1}}\frac{2}{\rho _1\rho _2}a(x,t,i\rho _1)+\sqrt{\frac{\rho _1}{\rho _2}}\frac{2}{\rho _2\rho _1}a(x,t,i\rho _2)\right].\hfill \end{array}$$
(2.26)
In this case, the even symmetry of the potential $`V_0(x,t)`$ guarantees that we may choose one basis element to be even and the other to be odd; we are using superscripts “(e)” and “(o)” to refer to even and odd functions of $`x`$ respectively. These bound state solutions of the linear Schrödinger equation are shown in Figures 3 and 4.
The two bound state modes are Bloch functions in $`t`$, with the same Floquet multiplier $`\mathrm{exp}(2i\beta _\mathrm{b}L)=\mathrm{exp}(2i\rho _1^2L)=\mathrm{exp}(2i\rho _2^2L)`$. Note that, in reference to the remark made at the end of §2.1, the function $`\mathrm{\Psi }_\mathrm{b}^{(\mathrm{e})}(x,t)`$ is proportional to $`\psi =2ib^{(1)}(x,t)`$, which is a two-soliton solution of the nonlinear Schrödinger equation
$$i_t\psi +\frac{1}{2}_x^2\psi +|\psi |^2\psi =0.$$
(2.27)
Correspondingly, $`V_0(x,t)=|\psi |^2`$ is the self-consistent potential.
It will also be useful to decompose the continuum into odd and even parts. Using the fact that $`a(x,t,\lambda )=a(x,t,\lambda )`$, define
$$\begin{array}{ccc}\hfill \mathrm{\Psi }_\mathrm{d}^{(\mathrm{e})}(x,t,\lambda )& & \frac{1}{\sqrt{2\pi (\lambda ^2+\rho _1^2)(\lambda ^2+\rho _2^2)}}(a(x,t,\lambda )+a(x,t,\lambda )),\hfill \\ & & \\ \hfill \mathrm{\Psi }_\mathrm{d}^{(\mathrm{o})}(x,t,\lambda )& & \frac{1}{\sqrt{2\pi (\lambda ^2+\rho _1^2)(\lambda ^2+\rho _2^2)}}(a(x,t,\lambda )a(x,t,\lambda )),\hfill \end{array}$$
(2.28)
where $`\lambda 0`$. The solutions $`\mathrm{\Psi }_\mathrm{d}^\mathrm{e}(x,t,\lambda )`$ and $`\mathrm{\Psi }_\mathrm{d}^{(\mathrm{o})}(x,t,\lambda )`$ are also Bloch functions in $`t`$ of period $`L`$ with Floquet multiplier $`\mathrm{exp}(2i\lambda ^2L)`$.
These solutions of the unperturbed problem have the following inner products :
$$\begin{array}{ccc}\hfill \mathrm{\Psi }_\mathrm{b}^{(\mathrm{e})}(,t),\mathrm{\Psi }_\mathrm{d}^{(\mathrm{e})}(,t,\lambda )& =& 0,\hfill \\ & & \\ \hfill \mathrm{\Psi }_\mathrm{b}^{(\mathrm{e})}(,t),\mathrm{\Psi }_\mathrm{b}^{(\mathrm{e})}(,t)& =& 1,\hfill \\ & & \\ \hfill \mathrm{\Psi }_\mathrm{d}^{(\mathrm{e})}(,t,\lambda ),\mathrm{\Psi }_\mathrm{d}^{(\mathrm{e})}(,t,\eta )& =& \delta (\lambda \eta ).\hfill \end{array}$$
(2.29)
In this latter relation it is assumed that both $`\lambda `$ and $`\eta `$ are positive real. Similar relations hold among the odd solutions, and of course everything even is orthogonal to everything odd. If the perturbation $`W(x,t)`$ is also even in $`x`$, then this observation will allow us to treat the even and odd parts of the field $`f(x,t)`$ in isolation to all orders in perturbation theory.
In our subsequent analysis of the coupled mode equations for this family of periodic potentials, we shall assume that the perturbation $`W(x,t)`$ is also an even function of $`x`$, and thus restrict attention to the subspace of initial conditions, $`f(x,0)`$ which are either even or odd in $`x`$. By the spatial symmetry of $`V=V_0+W`$, $`f(x,t)`$ has the same parity as $`f(x,0)`$. In analogy with the above derivation of coupled mode equations, we can then expand $`f(x,t)`$ in terms of (even or odd) modes of the unperturbed problem:
$$f(x,t)=B_\mathrm{b}^{(\alpha )}(t)\mathrm{\Psi }_\mathrm{b}^{(\alpha )}(x,t)+_0^{\mathrm{}}B_\mathrm{d}^{(\alpha )}(t,\lambda )\mathrm{\Psi }_\mathrm{d}^{(\alpha )}(x,t,\lambda )𝑑\lambda ,$$
(2.30)
where $`\alpha =\mathrm{e}`$ if $`f(x,0)`$ is even and $`\alpha =\mathrm{o}`$ if $`f(x,0)`$ is odd. Coupled mode equations for the amplitudes $`B_\mathrm{b}^{(\alpha )}(t),B_\mathrm{d}^{(\alpha )}(t,\lambda ),\lambda `$ analogous to those derived in the absence of any particular symmetry can then be derived by projecting the dynamical system (2.14) onto these even and odd basis modes. In §3 we show, by using the Floquet factorization of the unitary evolution associated with the unperturbed dynamics, that these coupled mode equations can be reexpressed as the following system (c.f. the system (3.15)) which is more amenable to our techniques:
$$\begin{array}{ccccc}\hfill i_tA_\mathrm{b}(t)& +& 2\beta _\mathrm{b}A_\mathrm{b}(t)& =& M(t)A_\mathrm{b}(t)+_0^{\mathrm{}}N^{(\mathrm{p})}(t,\lambda )A_\mathrm{d}(t,\lambda )𝑑\lambda ,\hfill \\ & & & & \\ \hfill i_tA_\mathrm{d}(t,\eta )& & 2\eta ^2A_\mathrm{d}(t,\eta )& =& N^{(\mathrm{p})}(t,\eta )^{}A_\mathrm{b}(t)+_0^{\mathrm{}}K^{(\mathrm{p})}(t,\eta ,\lambda )A_\mathrm{d}(t,\lambda )𝑑\lambda ,\hfill \end{array}$$
(2.31)
where
$$A_\mathrm{b}(t)B_\mathrm{b}(t)e^{2i\beta _\mathrm{b}t},A_\mathrm{d}(t,\lambda )B_\mathrm{d}(t,\lambda )e^{2i\lambda ^2t},$$
(2.32)
and where the scalar coefficients, all periodic with period $`L`$, are
$$\begin{array}{ccc}\hfill M(t)& =& \mathrm{\Psi }_\mathrm{b}(,t),W(,t)\mathrm{\Psi }_\mathrm{b}(,t)\hfill \\ & & \\ \hfill N^{(\mathrm{p})}(t,\lambda )& =& \mathrm{\Psi }_\mathrm{b}(,t),W(,t)\mathrm{\Psi }_\mathrm{d}(,t,\lambda )e^{2i(\lambda ^2+\beta _\mathrm{b})t}\hfill \\ & & \\ \hfill K^{(\mathrm{p})}(t,\eta ,\lambda )& =& \mathrm{\Psi }_\mathrm{d}(,t,\eta ),W(,t)\mathrm{\Psi }_\mathrm{d}(,t,\lambda )e^{2i(\lambda ^2\eta ^2)t}.\hfill \end{array}$$
(2.33)
Remark: To avoid cumbersome formulae, we have omitted the superscripts (o) and (e), with the understanding that the amplitudes correspond to either one type or the other, depending on the parity of $`f(x,0)`$. $`\mathrm{}`$
The system (2.31) may be viewed as that governing a family of oscillators: a single discrete oscillator whose amplitude is $`A_\mathrm{b}(t)`$ coupled to a continuum of oscillators with amplitudes $`A_\mathrm{d}(t,\eta ),\eta _+`$.
In §4 we shall analyze the coupled mode system (2.31), and determine the detailed asymptotic behavior of its solutions for small $`W(x,t)`$ over different timescales.
## 3 Coupled mode equations for periodic potentials.
Consider a dynamical system of the form (1.4), where both the unperturbed and the perturbed potential are time-periodic with the same period $`L`$. We will encounter a concrete example of such a problem in §5. Floquet theory suggests the introduction of a new time-periodic basis, with respect to which the problem (2.14) becomes a periodic perturbation of autonomous Hamiltonian system. This change of basis transforms the problem at hand into one similar to that treated in and . Similar methods are used along with resonance theory in a weakly nonlinear setting in .
### 3.1 Floquet factorization.
Let $`𝒰(t)`$ denote the unitary evolution operator (or propagator) of the unperturbed problem, so that for any $`L^2()`$ function $`f(x)`$, $`f(x,t)=𝒰(t)f(x)`$ is the solution of the unperturbed problem with $`f(x,0)=f(x)`$. As a consequence of the periodicity, the evolution operator can be factored into two operators on $`L^2()`$:
$$𝒰(t)=𝒫(t)e^{it}$$
(3.1)
where $`𝒫(t+L)=𝒫(t)`$, and $``$ is independent of $`t`$. This factorization can be motivated by the observation that by periodicity, there is an operator $``$ satisfying $`𝒰(t+L)=𝒰(t)`$, and that by setting $`t=0`$ in fact one has $`=𝒰(L)`$. Since $`𝒰(L)`$ is unitary, one can find a self-adjoint operator $``$ such that $`=𝒰(L)=e^{iL}`$. This operator $``$ in turn defines the abelian unitary group $`e^{it}`$. Now it is easy to see that $`𝒫(t)=𝒰(t)e^{it}`$ is a unitary operator satisfying $`𝒫(t+L)=𝒫(t)`$.
Let $`𝒲(t)`$ be the operator of multiplication by the correction to the potential $`W(x,t)`$, and set $`y(x,t)=𝒫(t)^{}f(x,t)`$. Then, the perturbed equation (1.4) becomes
$$i_tyy=\stackrel{~}{𝒲}(t)y,$$
(3.2)
where $`\stackrel{~}{𝒲}(t)𝒫(t)^{}𝒲(t)𝒫(t)`$, a “dressing” of $`𝒲(t)`$. The form of (3.2) is similar to the type of problem treated in and . The “unperturbed Hamiltonian” $``$ is time-independent, and self-adjoint. The perturbation $`\stackrel{~}{𝒲}(t)`$ is localized, self-adjoint, and time-periodic because the periods of $`𝒫(t)`$ and $`𝒲(t)`$ are equal. Typically the perturbation contains frequency components at all overtones of the fundamental frequency, and thus the version of the theory described in is most appropriate. The key qualitative difference between the present situation and that treated in is that here the unperturbed operator can have multiple bound states. We will soon introduce a symmetry that removes this difficulty from the scope of this paper. However, the methods of and can be extended to give results on radiation damping due to the coupling of multiple discrete modes to the continuum for a general class of spatially localized and time-dependent perturbations .
In fact, one can simplify the problem even further by invoking the spectral theorem for the self-adjoint operator $``$. This guarantees the existence of an isomorphism $`𝒱:L^2()L^2(\mathrm{\Sigma },d\mu )`$ to the space of square integrable functions on some set $`\mathrm{\Sigma }`$ with measure $`d\mu `$, such that $`𝒱=𝒯𝒱`$ where $`𝒯`$ is a real diagonal operator on $`L^2(\mathrm{\Sigma },d\mu )`$ (i.e. an operator of multiplication by a function from $`\mathrm{\Sigma }`$ to $``$). Setting $`z(t)=𝒱y(t)`$, we find the equation
$$i_tz𝒯z=𝒱\stackrel{~}{𝒲}(t)𝒱^{}z.$$
(3.3)
In quantum mechanics, making the transformation from (2.14) to (3.3) to facilitate the study of perturbations is known as going from the Schrödinger picture into the interaction picture.
In the particular example we will analyze in detail, arising from even perturbations of the two-soliton even potential described at the end of §2, the operator $``$ has a single degenerate eigenvalue of $`2\beta _\mathrm{b}<0`$ of geometric multiplicity two. By restricting separately to even and odd spaces of initial conditions (possible because the potential $`V_0(x,t)+W(x,t)`$ is symmetric in $`x`$) the problem is reduced to one which, formally, is precisely of the type studied in . We may then apply the methods developed in (subject to some appropriate hypotheses) without modification.
We now use our explicit knowledge developed in §2 of the unitary propagator $`𝒰(t)`$ corresponding to a time-periodic separable potential $`V_0(x,t)`$ to find the operators $`𝒫(t)`$ and $``$, and then to diagonalize $``$. This effectively implements the program described above and casts the perturbed problem (2.14) in a form (3.3) more suitable for analysis. We begin with the observation that each element of the basis of solutions of the unperturbed problem is a Bloch function or Floquet mode. We have:
$$\mathrm{\Psi }_\mathrm{d}(x,t+L,\lambda )=e^{2i\lambda ^2L}\mathrm{\Psi }_\mathrm{d}(x,t,\lambda )$$
(3.4)
where $`\lambda `$ is arbitrary real. Also, we have
$$a(x,t+L,i\rho _k)=e^{2i\rho _k^2L}a(x,t,i\rho _k).$$
(3.5)
Note that the commensurability relations (2.11) imply quite generally that the Floquet multipliers $`\mathrm{exp}(2i\rho _k^2L)`$ are all equal. This generalizes the observation made above in the context of the two-soliton potentials. This means that the entire $`M`$-dimensional subspace of bound states consists of degenerate Floquet modes. In particular, the elements of any orthonormal basis $`\{\mathrm{\Psi }_{\mathrm{b},k}(x,t),k=1,\mathrm{},M\}`$ have the same Floquet multiplier<sup>4</sup><sup>4</sup>4Recall that the Floquet exponents are not unique but that the Floquet multipliers are. Identification of the Floquet exponents with a single number $`\beta _\mathrm{b}>0`$ amounts to a particular choice of branch of the logarithm. $`\mathrm{exp}(2i\beta _\mathrm{b}L)`$. It now follows from (3.4) and (3.5) that the functions defined by
$$\begin{array}{ccc}\hfill \mathrm{\Psi }_\mathrm{d}^{(\mathrm{p})}(x,t,\lambda )& & e^{2i\lambda ^2t}\mathrm{\Psi }_\mathrm{d}(x,t,\lambda ),\hfill \\ & & \\ \hfill \mathrm{\Psi }_{\mathrm{b},k}^{(\mathrm{p})}(x,t)& & e^{2i\beta _\mathrm{b}t}\mathrm{\Psi }_{\mathrm{b},k}(x,t),\hfill \end{array}$$
(3.6)
are time-periodic with period $`L`$, as denoted by the superscript “(p)”.
As described in §2, the solution of the unperturbed problem with initial data $`f(x)`$ is expanded as
$$\begin{array}{ccc}\hfill f(x,t)& =& 𝒰(t)f(x)\hfill \\ & & \\ & =& \underset{k=1}{\overset{M}{}}\mathrm{\Psi }_{\mathrm{b},k}(,0),f()\mathrm{\Psi }_{\mathrm{b},k}(x,t)+_{\mathrm{}}^{\mathrm{}}\mathrm{\Psi }_\mathrm{d}(,0,\lambda ),f()\mathrm{\Psi }_\mathrm{d}(x,t,\lambda )𝑑\lambda \hfill \\ & & \\ & =& \underset{k=1}{\overset{M}{}}\mathrm{\Psi }_{\mathrm{b},k}^{(\mathrm{p})}(,0),e^{2i\beta _\mathrm{b}t}f()\mathrm{\Psi }_{\mathrm{b},k}^{(\mathrm{p})}(x,t)+_{\mathrm{}}^{\mathrm{}}\mathrm{\Psi }_\mathrm{d}^{(\mathrm{p})}(,0,\lambda ),e^{2i\lambda ^2t}f()\mathrm{\Psi }_\mathrm{d}^{(\mathrm{p})}(x,t,\lambda )𝑑\lambda .\hfill \end{array}$$
(3.7)
We now use the completeness relation at $`t=0`$ to factor $`𝒰(t)`$ as $`𝒫(t)e^{it}`$ where
$$e^{it}f(x)=\underset{k=1}{\overset{M}{}}\mathrm{\Psi }_{\mathrm{b},k}^{(\mathrm{p})}(,0),e^{2i\beta _\mathrm{b}t}f()\mathrm{\Psi }_{\mathrm{b},k}^{(\mathrm{p})}(x,0)+_{\mathrm{}}^{\mathrm{}}\mathrm{\Psi }_\mathrm{d}^{(\mathrm{p})}(,0,\lambda ),e^{2i\lambda ^2t}f()\mathrm{\Psi }_\mathrm{d}^{(\mathrm{p})}(x,0,\lambda )𝑑\lambda ,$$
(3.8)
and
$$𝒫(t)g(x)=\underset{k=1}{\overset{M}{}}\mathrm{\Psi }_{\mathrm{b},k}^{(\mathrm{p})}(,0),g()\mathrm{\Psi }_{\mathrm{b},k}^{(\mathrm{p})}(x,t)+_{\mathrm{}}^{\mathrm{}}\mathrm{\Psi }_\mathrm{d}^{(\mathrm{p})}(,0,\lambda ),g()\mathrm{\Psi }_\mathrm{d}^{(\mathrm{p})}(x,t,\lambda )𝑑\lambda .$$
(3.9)
We have used, several times, the fact that at $`t=0`$ there is no distinction between the basis elements and their periodic counterparts defined by (3.6). It is easy to see that $`f(x,t)=𝒰(t)f(x)=𝒫(t)e^{it}f(x)`$ is the solution of the unperturbed initial value problem with data $`f(x)L^2()`$ and that $`𝒫(t)`$ is periodic with period $`L`$ and $`𝒰(L)=e^{iL}`$.
The generator of the abelian unitary group $`e^{it}`$ is
$$\begin{array}{ccc}\hfill f(x)& =& i\frac{d}{dt}e^{it}f(x)|_{t=0}\hfill \\ & & \\ & =& \underset{k=1}{\overset{M}{}}\mathrm{\Psi }_{\mathrm{b},k}^{(\mathrm{p})}(,0),2\beta _\mathrm{b}f()\mathrm{\Psi }_{\mathrm{b},k}^{(\mathrm{p})}(x,0)+_{\mathrm{}}^{\mathrm{}}\mathrm{\Psi }_\mathrm{d}^{(\mathrm{p})}(,0,\lambda ),2\lambda ^2f()\mathrm{\Psi }_\mathrm{d}^{(\mathrm{p})}(x,0,\lambda )𝑑\lambda \hfill \\ & & \\ & =& \underset{k=1}{\overset{M}{}}\mathrm{\Psi }_{\mathrm{b},k}(,0),2\beta _\mathrm{b}f()\mathrm{\Psi }_{\mathrm{b},k}(x,0)+_{\mathrm{}}^{\mathrm{}}\mathrm{\Psi }_\mathrm{d}(,0,\lambda ),2\lambda ^2f()\mathrm{\Psi }_\mathrm{d}(x,0,\lambda )𝑑\lambda .\hfill \end{array}$$
(3.10)
In the last step we have dropped the superscripts “(p)” since everything is evaluated at $`t=0`$. This formula for the self-adjoint operator $``$ makes clear its spectral decomposition. The isomorphism $`𝒱`$ takes a function $`g(x)L^2()`$ to a function $`A_\mathrm{d}(\lambda )`$ for $`\lambda `$ and a set of $`M`$ numbers $`A_{\mathrm{b},k}`$ for $`k=1,\mathrm{},M`$ defined by
$$A_\mathrm{d}(\lambda )\mathrm{\Psi }_\mathrm{d}^{(\mathrm{p})}(,0,\lambda ),g(),$$
(3.11)
and for $`k=1,\mathrm{},M`$,
$$A_{\mathrm{b},k}\mathrm{\Psi }_{\mathrm{b},k}^{(\mathrm{p})}(,0),g().$$
(3.12)
The diagonal operator $`𝒯`$ is then simply defined by
$$𝒯\left[\begin{array}{c}A_\mathrm{d}(\lambda )\\ A_{\mathrm{b},1}\\ \mathrm{}\\ A_{\mathrm{b},M}\end{array}\right]=\left[\begin{array}{cccc}2\lambda ^2& & & \\ & 2\beta _\mathrm{b}& & \\ & & \mathrm{}& \\ & & & 2\beta _\mathrm{b}\end{array}\right]\left[\begin{array}{c}A_\mathrm{d}(\lambda )\\ A_{\mathrm{b},1}\\ \mathrm{}\\ A_{\mathrm{b},M}\end{array}\right].$$
(3.13)
It is now easy to use the definition of the unitary periodic operator $`𝒫(t)`$ and the unitary isomorphism $`𝒱`$, along with the completeness relation to compute the dressed operator $`𝒱\stackrel{~}{𝒲}(t)𝒱^{}`$ and thus write the perturbed problem (2.14) in the simple form (3.3). The dynamical unknowns are in the range of $`𝒱`$, the space $`L^2(\mathrm{\Sigma },d\mu )`$, and are given in terms of $`f(x,t)`$, the solution of (2.14), by
$$\begin{array}{ccccc}\hfill A_{\mathrm{b},k}& =& (𝒱𝒫(t)^{}f(,t))_{\mathrm{b},k}& =& \mathrm{\Psi }_{\mathrm{b},k}^{(\mathrm{p})}(,t),f(,t),\hfill \\ \hfill A_\mathrm{d}(\lambda )& =& (𝒱𝒫(t)^{}f(,t))_\mathrm{d}(\lambda )& =& \mathrm{\Psi }_\mathrm{d}^{(\mathrm{p})}(,t,\lambda ),f(,t).\hfill \end{array}$$
(3.14)
When $`f(x,t)`$ satisfies (2.14), these quantities satisfy the system
$$\begin{array}{ccc}\hfill i_t\stackrel{}{A}_\mathrm{b}+2\beta _\mathrm{b}\stackrel{}{A}_\mathrm{b}& =& 𝐌(t)\stackrel{}{A}_\mathrm{b}+_{\mathrm{}}^{\mathrm{}}A_\mathrm{d}(\lambda )\stackrel{}{N}^{(\mathrm{p})}(t,\lambda )𝑑\lambda ,\hfill \\ & & \\ \hfill i_tA_\mathrm{d}(\eta )2\eta ^2A_\mathrm{d}(\eta )& =& \stackrel{}{N}^{(\mathrm{p})}(t,\eta )^{}\stackrel{}{A}_\mathrm{b}+_{\mathrm{}}^{\mathrm{}}K^{(\mathrm{p})}(t,\eta ,\lambda )A_\mathrm{d}(\lambda )𝑑\lambda ,\hfill \end{array}$$
(3.15)
where $`\stackrel{}{A}_\mathrm{b}`$ is the vector of components $`A_{\mathrm{b},1},\mathrm{},A_{\mathrm{b},M}`$, and the time-periodic matrix elements are defined in terms of (2.17) by
$$\begin{array}{ccc}\hfill \stackrel{}{N}^{(\mathrm{p})}(t,\lambda )& & e^{2i(\lambda ^2+\beta _\mathrm{b})t}\stackrel{}{N}(t,\lambda ),\hfill \\ & & \\ \hfill K^{(\mathrm{p})}(t,\eta ,\lambda )& & e^{2i(\lambda ^2\eta ^2)t}K(t,\eta ,\lambda ).\hfill \end{array}$$
(3.16)
The periodicity of these matrix elements when $`W(x,t)`$ is periodic with period $`L`$ is also clear from these explicit formulas and the Bloch relations (3.4) and (3.5) for the basis of solutions; these imply similar ones for the matrix elements defined by (2.17). We have
$$\begin{array}{ccc}\hfill 𝐌(t+L)& =& 𝐌(t),\hfill \\ & & \\ \hfill \stackrel{}{N}(t+L,\lambda )& =& e^{2i(\lambda ^2+\beta _\mathrm{b})L}\stackrel{}{N}(t,\lambda ),\hfill \\ & & \\ \hfill K(t+L,\eta ,\lambda )& =& e^{2i(\eta ^2\lambda ^2)L}K(t,\eta ,\lambda ).\hfill \end{array}$$
(3.17)
Of course, the right-hand side of (3.15) is just the operator $`𝒱\stackrel{~}{𝒲}(t)𝒱^{}`$ operating on the dynamical unknowns. Similarly, the perturbation operator $`\stackrel{~}{𝒲}(t)`$ operating in the space $`L^2()`$ can be explicitly written as
$$\begin{array}{ccc}\hfill \stackrel{~}{𝒲}(t)f(x)& =& \underset{k=1}{\overset{M}{}}\underset{l=1}{\overset{M}{}}M_{k,l}(t)\mathrm{\Psi }_{\mathrm{b},l}(,0),f()\mathrm{\Psi }_{\mathrm{b},k}(x,0)\hfill \\ & & \\ & +& \underset{k=1}{\overset{M}{}}_{\mathrm{}}^{\mathrm{}}N_k^{(\mathrm{p})}(t,\eta )\mathrm{\Psi }_\mathrm{d}(,0,\eta ),f()\mathrm{\Psi }_{\mathrm{b},k}(x,0)𝑑\eta \hfill \\ & & \\ & +& \underset{l=1}{\overset{M}{}}_{\mathrm{}}^{\mathrm{}}N_l^{(\mathrm{p})}(t,\lambda )^{}\mathrm{\Psi }_{\mathrm{b},l}(,0),f()\mathrm{\Psi }_\mathrm{d}(x,0,\lambda )𝑑\lambda \hfill \\ & & \\ & +& _{\mathrm{}}^{\mathrm{}}_{\mathrm{}}^{\mathrm{}}K^{(\mathrm{p})}(t,\lambda ,\eta )\mathrm{\Psi }_\mathrm{d}(,0,\eta ),f()\mathrm{\Psi }_\mathrm{d}(x,0,\lambda )𝑑\lambda 𝑑\eta .\hfill \end{array}$$
(3.18)
For the special choice of $`V_0(x,t)`$ discussed at the end of §2, evenness implies that there is one bound state of each parity. If the perturbation $`W(x,t)`$ also has even symmetry in $`x`$, the coupled mode system (3.15) reduces to a system of the type (2.31) when the initial condition is restricted to either even or odd parity. It is easily checked that the unknowns as defined above correspond exactly to those defined in §2 for the system (2.31).
## 4 Analysis of the Coupled Mode Equations
In this section we study the structural instability of the even and odd breather modes introduced at the end of §2 associated with the two-soliton time-periodic even potentials. We first give a simple argument valid for short times that in the presence of a perturbation $`W(x,t)`$ to the potential $`V_0(x,t)`$, the bound state begins to decay initially. We then seek to capture the dynamics for longer times, primarily to show that this initial phase of decay does not reverse itself, but takes on a different, exponentially decaying, character. The decay will be first calculated formally, using asymptotic expansions and the method of multiple scales. Then, using the results of Kirr and Weinstein , we show that at least in the odd case, it is possible to make statements about the decay process that are valid globally in time. In particular, these arguments will rigorously justify the formal results for the odd case, and will show that the exponential decay model is only a valid approximation until it becomes smaller than the dispersive part of the solution. The bound state ultimately dies algebraically in time, qualitatively indistinguishable from the dispersive components of the solution to which it is orthogonal.
### 4.1 Small time analysis. The watched pot effect.
A simple calculation carried out at the level of the coupled mode equations (2.31) shows that the effect of the perturbation is to cause the bound state to decay immediately both forward and backward in time. More complicated calculations will be required to show that the decay does not stop or reverse for longer times, although it takes on a different character. The approach in the small time analysis is simply to expand the solution in Taylor series:
$$\begin{array}{ccc}\hfill A_\mathrm{b}(t)& =& A_\mathrm{b}(0)+c_1t+c_2t^2+O(t^3),\hfill \\ \hfill A_\mathrm{d}(\lambda ,t)& =& d_1(\lambda )t+O(t^2),\hfill \end{array}$$
(4.1)
and use the (known) Taylor expansions of the matrix elements, in particular,
$$\begin{array}{ccc}\hfill M(t)& =& M(0)+M^{}(0)t+O(t^2),\hfill \\ \hfill N^{(\mathrm{p})}(t,\lambda )& =& N^{(\mathrm{p})}(0,\lambda )+O(t).\hfill \end{array}$$
(4.2)
Substituting these series into (2.31), one finds:
$$\begin{array}{c}\left[ic_1+2\beta _\mathrm{b}A(0)M(0)A_\mathrm{b}(0)\right]+\hfill \\ \\ \left[2ic_2+2\beta _\mathrm{b}c_1M(0)c_1M^{}(0)A_\mathrm{b}(0)_0^{\mathrm{}}d_1(\lambda )N^{(\mathrm{p})}(0,\lambda )𝑑\lambda \right]t=O(t^2),\hfill \\ \\ id_1(\eta )N^{(\mathrm{p})}(0,\eta )^{}A_\mathrm{b}(0)=O(t).\hfill \end{array}$$
(4.3)
Solving for $`c_1`$ and $`c_2`$ yields an approximation for $`A_\mathrm{b}(t)`$, valid for small $`t`$:
$$\begin{array}{c}A_\mathrm{b}(t)=A_\mathrm{b}(0)[1i(M(0)2\beta _\mathrm{b})t\hfill \\ \\ \frac{1}{2}(iM^{}(0)+(M(0)2\beta _\mathrm{b})^2+_0^{\mathrm{}}|N^{(\mathrm{p})}(0,\lambda )|^2d\lambda )t^2+O(t^3)].\hfill \end{array}$$
(4.4)
It easily follows that
$$|A_\mathrm{b}(t)|^2=|A_\mathrm{b}(0)|^2\left[1t^2_0^{\mathrm{}}|N^{(\mathrm{p})}(0,\lambda )|^2𝑑\lambda +O(t^3)\right].$$
(4.5)
Note that smallness of the perturbation is not exploited in these calculations. This Taylor expansion shows that the initial phase of the evolution is a process of radiative decay, since $`|A_\mathrm{b}(t)|^2<|A_\mathrm{b}(0)|^2`$ for all nonzero $`t`$ in some neighborhood of $`t=0`$. The decay is symmetric in time.
The fact that the decay is an order $`O(t^2)`$ effect is quite general<sup>5</sup><sup>5</sup>5In the general setting, the decay is a simple consequence of the Cauchy-Schwarz inequality. One supposes that $`𝒰(t)`$ is the unitary propagator of the possibly time-dependent unperturbed problem:
$$i𝒰_t(t)\varphi ^0=_0(t)𝒰(t)\varphi ^0,$$
for all states $`\varphi ^0`$. One then considers the perturbed equation
$$i\psi _t=(_0(t)+𝒲(t))\psi ,$$
by setting $`\psi (t)=𝒰(t)\varphi (t)`$, giving the “interaction picture” equation
$$i\varphi _t=𝒰(t)^{}𝒲(t)𝒰(t)\varphi ,$$
which one solves by Taylor series in $`t`$. The result is:
$$\varphi (t)=\left(i𝒲(0)t+\frac{t^2}{2}\left(i𝒲^{}(0)𝒲(0)^2+[_0(0),𝒲(0)]\right)+O(t^3)\right)\varphi (0).$$
The probability of remaining in the unperturbed state is then found to be (using self-adjointness of both $`𝒲(0)`$ and $`_0(0)`$)
$$\begin{array}{ccc}\hfill |𝒰(t)\varphi (0),𝒰(t)\varphi (t)|^2& =& |\varphi (0),\varphi (t)|^2\hfill \\ & =& \varphi (0)_2^4(𝒲(0)\varphi (0)_2^2\varphi (0)_2^2|\varphi (0),𝒲(0)\varphi (0)|^2)t^2+O(t^3).\hfill \end{array}$$
This quantity is initially decreasing in time as a consequence of the Cauchy-Schwarz inequality. and well-known in the perturbation theory of stationary Schrödinger equations. It has an interesting interpretation in the quantum theory of ideal measurements, the so-called “watched pot effect”. Suppose that an ideal measurement is made at some point during the evolution of the wave function to determine whether the state is bound, and the measurement yields a positive result. The probability of a positive result at time $`t`$ is $`|A_\mathrm{b}(t)|^2/|A_\mathrm{b}(0)|^2`$. The theory of ideal measurements says that as a consequence of the measurement disturbing the system, the wave function “collapses” upon a positive result to the bound state, and evolution of the wave function according to the Schrödinger equation continues from this “reset” bound state. One may then try to determine the asymptotic effect of making many such measurements in a finite time interval. In particular, we can ask about the limiting probability of finding the system in the bound state after each of $`n`$ ideal measurements performed at times $`t_n=T/n`$, as $`n\mathrm{}`$. After each positive result, the wave function collapses and the experiment is restarted. The Schrödinger evolution takes place over short time intervals so it is appropriate to replace the probability in each interval $`p(t)`$ by its short-time approximation $`p(t)=1(\alpha t)^2+O(t^3)`$. The $`n`$ measurements are independent events, so the probability of always finding the system bound after each measurement is simply
$$P_n=p(T/n)^n.$$
(4.6)
Because the “time slice” decay probability $`1p(t)`$ is quadratic in $`t`$, $`P_n`$ tends to unity<sup>6</sup><sup>6</sup>6The superlinear nature of the decay probability is important. If $`p(t)=1|\alpha t|+O(t^2)`$, then $`P_n`$ tends to $`e^{|\alpha T|}<1`$ instead. as $`n\mathrm{}`$, regardless of the value of $`T`$. So if the measurements are performed infinitely often, the decay of the bound state never occurs. The quantum “watched pot” never boils.
### 4.2 Multiple scales analysis.
We begin the multiple scales analysis by assuming that the correction $`W(x,t)=W(x,t;ϵ)`$ to the potential energy has an expansion in a small parameter, $`ϵ`$ (see, for example, equation (5.11) ):
$$W(x,t;ϵ)=ϵW_1(x,t)+O(ϵ^2).$$
It then follows that the coupling coefficient functions in (2.31) have formal expansions for small $`ϵ`$:
$$\begin{array}{ccc}\hfill M(t)& =& ϵM_1(t)+ϵ^2M_2(t)+O(ϵ^3),\hfill \\ \hfill N^{(\mathrm{p})}(t,\lambda )& =& ϵN_1^{(\mathrm{p})}(t,\lambda )+O(ϵ^2),\hfill \\ \hfill K^{(\mathrm{p})}(t,\eta ,\lambda )& =& ϵK_1^{(\mathrm{p})}(t,\eta ,\lambda )+O(ϵ^2).\hfill \end{array}$$
(4.7)
Here, $`M_1(t),N_1^{(\mathrm{p})}(t)`$ and $`K_1^{(\mathrm{p})}(t,\eta ,\lambda )`$ correspond to the expressions for $`M(t)`$, $`N^{(\mathrm{p})}`$ and $`K^{(\mathrm{p})}`$ in (2.33) with $`W`$ replaced by $`W_1`$.
The amplitudes $`A_\mathrm{b}(t;ϵ)`$ and $`A_\mathrm{d}(t,\lambda ;ϵ)`$ are assumed to have asymptotic expansions of the form
$$\begin{array}{ccc}\hfill A_\mathrm{b}(t;ϵ)& =& A_\mathrm{b}^{(0)}(T_0,T_1,T_2,\mathrm{})+ϵA_\mathrm{b}^{(1)}(T_0,T_1,T_2,\mathrm{})+ϵ^2A_\mathrm{b}^{(2)}(T_0,T_1,T_2,\mathrm{})+O(ϵ^3),\hfill \\ \hfill A_\mathrm{d}(t,\lambda ;ϵ)& =& A_\mathrm{d}^{(0)}(T_0,T_1,T_2,\mathrm{},\lambda )+ϵA_\mathrm{d}^{(1)}(T_0,T_1,T_2,\mathrm{},\lambda )+O(ϵ^2),\hfill \end{array}$$
(4.8)
where the $`T_kϵ^kt`$ are time scale variables. Such expansions of given functions $`A_\mathrm{b}(t;ϵ)`$ and $`A_\mathrm{d}(t,\lambda ;ϵ)`$ are highly nonunique. However the guiding principle of the method of multiple scales (see, for example, ) stipulates that the dependence of the various terms on the “slow” times $`T_1`$, $`T_2`$, and so on is chosen so that each term is uniformly bounded as a function of the “fast” time $`T_0`$. This procedure is quite systematic, and is supposed to keep the error terms in any truncation uniformly small in time intervals where $`T_k`$ is bounded for some $`k`$ as $`ϵ`$ tends to zero. We will see by comparison with the rigorous results that this formal procedure indeed works as advertised.
One now substitutes these expansions into the system (2.31) and expands the time derivative operating on the expansion coefficients in (4.8) according to the chain rule:
$$_t=_{T_0}+ϵ_{T_1}+ϵ^2_{T_2}+\mathrm{}.$$
(4.9)
The coupling coefficients, being all periodic functions of $`t`$ with period $`L`$ independent of $`ϵ`$, are taken to be explicit functions of $`t=T_0`$ only. Substituting them into (2.31) along with the expansions (4.8) and the chain rule formula (4.9), and equating terms with the same powers of $`ϵ`$ leads to a hierarchy of equations:
$`O(1):`$ $`\{\begin{array}{ccccc}\hfill i_{T_0}A_\mathrm{b}^{(0)}& +& 2\beta _\mathrm{b}A_\mathrm{b}^{(0)}& =& 0,\hfill \\ & & & & \\ \hfill i_{T_0}A_\mathrm{d}^{(0)}(\eta )& & 2\eta ^2A_\mathrm{d}^{(0)}(\eta )& =& 0,\hfill \end{array}`$ (4.13)
$`O(ϵ):`$ $`\{\begin{array}{ccccc}\hfill i_{T_0}A_\mathrm{b}^{(1)}& +& 2\beta _\mathrm{b}A_\mathrm{b}^{(1)}& =& i_{T_1}A_\mathrm{b}^{(0)}+M_1(T_0)A_\mathrm{b}^{(0)}\hfill \\ & & & & \\ & & & & +{\displaystyle _0^{\mathrm{}}}N_1^{(\mathrm{p})}(T_0,\lambda )A_\mathrm{d}^{(0)}(\lambda )𝑑\lambda ,\hfill \\ & & & & \\ \hfill i_{T_0}A_\mathrm{d}^{(1)}(\eta )& & 2\eta ^2A_\mathrm{d}^{(1)}(\eta )& =& i_{T_1}A_\mathrm{d}^{(0)}(\eta )+N_1^{(\mathrm{p})}(T_0,\eta )^{}A_\mathrm{b}^{(0)}\hfill \\ & & & & \\ & & & & +{\displaystyle _0^{\mathrm{}}}K_1^{(\mathrm{p})}(T_0,\eta ,\lambda )A_\mathrm{d}^{(0)}(\lambda )𝑑\lambda ,\hfill \end{array}`$ (4.21)
$`O(ϵ^2):`$ $`\{\begin{array}{ccccc}\hfill i_{T_0}A_\mathrm{b}^{(2)}& +& 2\beta _\mathrm{b}A_\mathrm{b}^{(2)}& =& i_{T_1}A_\mathrm{b}^{(1)}i_{T_2}A_\mathrm{b}^{(0)}\hfill \\ & & & & \\ & & & & +M_2(T_0)A_\mathrm{b}^{(0)}+M_1(T_0)A_\mathrm{b}^{(1)}\hfill \\ & & & & \\ & & & & +{\displaystyle _0^{\mathrm{}}}N_1^{(\mathrm{p})}(T_0,\lambda )A_\mathrm{d}^{(1)}(\lambda )𝑑\lambda \hfill \\ & & & & \\ & & & & +{\displaystyle _0^{\mathrm{}}}N_2^{(\mathrm{p})}(T_0,\lambda )A_\mathrm{d}^{(0)}(\lambda )𝑑\lambda ,\hfill \\ & & & & \\ \hfill i_{T_0}A_\mathrm{d}^{(2)}(\eta )& & 2\eta ^2A_\mathrm{d}^{(2)}(\eta )& =& \mathrm{},\hfill \end{array}`$ (4.31)
and so on. Our initial conditions are encoded in the expansions (4.8) as $`A_\mathrm{b}^{(0)}(0,0,0,\mathrm{})=A_{\mathrm{b0}}`$, $`A_\mathrm{d}^{(0)}(0,0,0,\mathrm{},\lambda )=0`$, and for $`j1`$, $`A_\mathrm{b}^{(j)}(0,0,0,\mathrm{})=A_\mathrm{d}^{(j)}(0,0,0,\mathrm{},\lambda )=0`$. We now proceed to solve the hierarchy sequentially.
Solving the equations (4.13) at order $`O(1)`$ subject to the initial conditions gives
$$A_\mathrm{b}^{(0)}=Ce^{2i\beta _\mathrm{b}T_0},A_\mathrm{d}^{(0)}(\eta )=0,$$
(4.32)
where $`C=C(T_1,T_2,\mathrm{})`$ satisfies the initial condition $`C(0,0,\mathrm{})=A_{\mathrm{b0}}`$ but is otherwise undetermined at this stage.
In the first of the two equations (4.21) appearing at $`O(ϵ)`$, it is natural to make the substitution
$$A_\mathrm{b}^{(1)}=fe^{2i\beta _\mathrm{b}T_0},$$
(4.33)
which leads to the equation
$$_{T_0}f=_{T_1}CiM_1(T_0)C.$$
(4.34)
Integrating with the use of the initial condition $`f(T_0=0)=0`$, and keeping in mind that $`T_1`$ and $`T_0`$ are to be thought of as independent variables, leads to the expression for $`A_\mathrm{b}^{(1)}`$:
$$A_\mathrm{b}^{(1)}=\left(_{T_1}CT_0iC_0^{T_0}M_1(s)𝑑s\right)e^{2i\beta _\mathrm{b}T_0}.$$
(4.35)
We need for this correction to be bounded as a function of $`T_0`$ so that the asymptotic expansion will be well-ordered for long times. Since $`M_1(s)`$ is a periodic function of period $`L`$, this requirement uniquely determines $`_{T_1}C`$:
$$_{T_1}C=i\overline{M_1}C,\overline{M_1}\frac{1}{L}_0^LM_1(s)𝑑s.$$
(4.36)
Thus,
$$C=De^{i\overline{M_1}T_1},D=D(T_2,\mathrm{}),D(0,\mathrm{})=A_{\mathrm{b0}}.$$
(4.37)
Putting together what we have for the bound state amplitude at this time,
$$A_\mathrm{b}^{(0)}=De^{i\overline{M_1}T_1}e^{2i\beta _\mathrm{b}T_0},A_\mathrm{b}^{(1)}=iDe^{i\overline{M_1}T_1}e^{2i\beta _\mathrm{b}T_0}_0^{T_0}\left(M_1(s)\overline{M_1}\right)𝑑s.$$
(4.38)
This has been the first application in our calculation of the guiding principle of the method of multiple scales, that dependence of expansion terms on “slow” times is chosen to ensure that the expansion terms are uniformly bounded with respect to the “fast” time $`T_0`$. Now we solve for the correction to the dispersive mode amplitude at this order (in fact the leading term) using the second of the equations (4.21). Substituting the expressions from the previous order and using the initial conditions gives a unique expression:
$$A_\mathrm{d}^{(1)}(\eta )=iDe^{i\overline{M_1}T_1}_0^{T_0}N_1^{(\mathrm{p})}(s,\eta )^{}e^{2i\beta _\mathrm{b}s}e^{2i\eta ^2(T_0s)}𝑑s.$$
(4.39)
Continuing systematically with the equation (4.31) for the bound state amplitude correction at order $`O(ϵ^2)`$, we substitute all the expressions known thus far and observe the utility of the change of variables
$$A_\mathrm{b}^{(2)}=he^{i\overline{M_1}T_1}e^{2i\beta _\mathrm{b}T_0}.$$
(4.40)
We find for $`h`$ the simple equation
$$_{T_0}h=D\left(M_1(T_0)\overline{M_1}\right)_0^{T_0}\left(M_1(s)\overline{M_1}\right)𝑑s_{T_2}DD\gamma (T_0),$$
(4.41)
where
$$\gamma (T_0)iM_2(T_0)+_0^{\mathrm{}}N_1^{(\mathrm{p})}(T_0,\lambda )_0^{T_0}e^{2i(\lambda ^2+\beta _\mathrm{b})(T_0s)}N_1^{(\mathrm{p})}(s,\lambda )^{}𝑑s𝑑\lambda .$$
(4.42)
Equation (4.41) can be analyzed as follows. By linearity, we can express $`h`$ as a sum: $`h=h_1+h_2`$, where
$$\begin{array}{ccc}\hfill _{T_0}h_1& =& D\left(M_1(T_0)\overline{M_1}\right)_0^{T_0}\left(M_1(s)\overline{M_1}\right)𝑑s,\hfill \\ & & \\ \hfill _{T_0}h_2& =& _{T_2}DD\gamma (T_0),\hfill \end{array}$$
(4.43)
and where we assume the initial conditions $`h_1(T_0=0)=h_2(T_0=0)=0`$. Integrating the equation for $`h_1`$ exactly using the initial condition gives:
$$h_1=\frac{D}{2}\left(_0^{T_0}(M_1(s)\overline{M_1})𝑑s\right)^2,$$
(4.44)
which is periodic, and in particular bounded, by periodicity of $`M_1(T_0)`$.
We now want to select the dependence of $`D`$ on the slow time $`T_2`$ such that $`h_2`$, as found from the second of equations (4.43), is a bounded function of $`T_0`$. Using the initial condition to integrate the equation for $`h_2`$ with respect to $`T_0`$ while holding $`T_2`$ fixed gives
$$h_2=T_0_{T_2}DD_0^{T_0}\gamma (s)𝑑s.$$
(4.45)
Clearly, the possibility of choosing $`D(T_2)`$ so that the expression (4.45) is bounded in $`T_0`$ depends on the behavior of $`\gamma (T_0)`$ in the limits $`T_0\pm \mathrm{}`$. We now study $`\gamma (T_0)`$ for large $`|T_0|`$.
We can compute the $`s`$-integral in (4.42) exactly if we introduce the Fourier series for the periodic function $`N_1^{(\mathrm{p})}(T_0,\lambda )`$:
$$N_1^{(\mathrm{p})}(T_0,\lambda )=\underset{k=\mathrm{}}{\overset{\mathrm{}}{}}N_{1,k}(\lambda )e^{2\pi ikT_0/L}.$$
(4.46)
Note that in terms of the Fourier coefficients of $`N^{(\mathrm{p})}(T_0,\lambda ;ϵ)`$ itself (see (4.7)), we have
$$N_{1,k}(\lambda )=\underset{ϵ0}{lim}ϵ^1N_k(\lambda ;ϵ).$$
(4.47)
Substituting the Fourier series into (4.42), integrating term by term with respect to $`s`$, and changing variables to $`\sigma =\lambda ^2`$, we arrive at
$$\gamma (T_0)=iM_2(T_0)+\underset{n,k=\mathrm{}}{\overset{\mathrm{}}{}}i_0^{\mathrm{}}\frac{N_{1,n}(\sqrt{\sigma })N_{1,k}(\sqrt{\sigma })^{}}{4\sqrt{\sigma }(\sigma \sigma _k)}\left[e^{2i(\sigma \sigma _n)T_0}e^{2\pi i(nk)T_0/L}\right]𝑑\sigma ,$$
(4.48)
where the resonances $`\sigma _n`$ are defined by
$$\sigma _n\pi n/L\beta _\mathrm{b}.$$
(4.49)
Note that for all terms having $`\sigma _k>0`$, the difference of the exponentials in the integrand vanishes for $`\sigma =\sigma _k`$, so there is no singularity. Moreover, the Fourier coefficients $`N_{1,n}(\lambda )`$ are by construction analytic functions of $`\lambda `$ for $`\lambda `$ in a sector including the real axis, and so the quantities $`N_{1,n}(\sqrt{\sigma })`$ are analytic in a neighborhood of the positive real $`\sigma `$ axis. This property extends to the whole integrand, and we may therefore deform the integration contour away from the real axis in an effort to study the behavior for large $`|T_0|`$ by a steepest descents type argument.
For positive $`T_0`$, we deform the contour into the lower half plane. For $`\delta >0`$, let $`C_+^\delta `$ be the contour consisting of the diagonal segment from $`0`$ to $`(1i)\delta `$ followed by the horizontal ray from $`(1i)\delta `$ to $`i\delta +\mathrm{}`$ (see Figure 5).
We have
$$\begin{array}{ccc}\hfill \gamma (T_0)& =& iM_2(T_0)+\underset{n,k=\mathrm{}}{\overset{\mathrm{}}{}}i_{C_+^\delta }\frac{N_{1,n}(\sqrt{\sigma })N_{1,k}(\sqrt{\sigma })^{}}{4\sqrt{\sigma }(\sigma \sigma _k)}\left[e^{2i(\sigma \sigma _n)T_0}e^{2\pi i(nk)T_0/L}\right]𝑑\sigma \hfill \\ & & \\ & =& iM_2(T_0)+\underset{n,k=\mathrm{}}{\overset{\mathrm{}}{}}i_{C_+^\delta }\frac{N_{1,n}(\sqrt{\sigma })N_{1,k}(\sqrt{\sigma })^{}}{4\sqrt{\sigma }(\sigma \sigma _k)}e^{2i(\sigma \sigma _n)T_0}𝑑\sigma \hfill \\ & & \\ & & \underset{n,k=\mathrm{}}{\overset{\mathrm{}}{}}i_{C_+^\delta }\frac{N_{1,n}(\sqrt{\sigma })N_{1,k}(\sqrt{\sigma })^{}}{4\sqrt{\sigma }(\sigma \sigma _k)}e^{2\pi i(nk)T_0/L}𝑑\sigma \hfill \\ & & \\ & =& \gamma _0(T_0)+\gamma _1^+(T_0)+\gamma _2^+(T_0),\hfill \end{array}$$
(4.50)
so that on the new contour $`C_+^\delta `$ the two integrals converge independently. The term $`\gamma _0(T_0)`$ is periodic in $`T_0`$ with period $`L`$ and mean value $`i\overline{M_2}`$. The term $`\gamma _2^+(T_0)`$ is also a periodic function of $`T_0`$ with period $`L`$. Its mean value is given by the terms in the sum with $`n=k`$:
$$\overline{\gamma _2^+}=\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}i_{C_+^\delta }\frac{|N_{1,n}(\sqrt{\sigma })|^2}{4\sqrt{\sigma }(\sigma \sigma _n)}𝑑\sigma .$$
(4.51)
Letting $`\delta `$ tend to zero does not alter the value of the integral, and then we may use the Plemelj-Sokhotski formula<sup>7</sup><sup>7</sup>7This is merely the distributional identity
$$(x\pm i0)^1=\mathrm{P}.\mathrm{V}.x^1i\pi \delta (x).$$
to evaluate the terms with $`\sigma _n>0`$ to find
$$\overline{\gamma _0+\gamma _2^+}=i\overline{M_2}i\mathrm{\Lambda }_2+\mathrm{\Gamma }_2,$$
(4.52)
where
$$\mathrm{\Lambda }_2\underset{n=\mathrm{}}{\overset{n_01}{}}_0^{\mathrm{}}\frac{|N_{1,n}(\sqrt{\sigma })|^2d\sigma }{4\sqrt{\sigma }(\sigma \sigma _n)}+\underset{n=n_0}{\overset{\mathrm{}}{}}\mathrm{P}.\mathrm{V}._0^{\mathrm{}}\frac{|N_{1,n}(\sqrt{\sigma })|^2d\sigma }{4\sqrt{\sigma }(\sigma \sigma _n)},$$
(4.53)
and
$$\mathrm{\Gamma }_2\frac{\pi }{4}\underset{n=n_0}{\overset{\mathrm{}}{}}\frac{|N_{1,n}(\sqrt{\sigma _n})|^2}{\sqrt{\sigma _n}}.$$
(4.54)
Finally, consider the term $`\gamma _1^+(T_0)`$. Its time integral, calculated term by term, is
$$\begin{array}{ccc}\hfill _0^{T_0}\gamma _1^+(s)𝑑s& =& \underset{n,k=\mathrm{}}{\overset{\mathrm{}}{}}_{C_+^\delta }\frac{N_{1,n}(\sqrt{\sigma })N_{1,k}(\sqrt{\sigma })^{}}{8\sqrt{\sigma }(\sigma \sigma _k)(\sigma \sigma _n)}\left[e^{2i(\sigma \sigma _n)T_0}1\right]𝑑\sigma \hfill \\ & & \\ & =& \underset{n,k=\mathrm{}}{\overset{\mathrm{}}{}}_{C_+^\delta }\frac{N_{1,n}(\sqrt{\sigma })N_{1,k}(\sqrt{\sigma })^{}}{8\sqrt{\sigma }(\sigma \sigma _k)(\sigma \sigma _n)}𝑑\sigma \hfill \\ & & \\ & & \underset{n,k=\mathrm{}}{\overset{\mathrm{}}{}}_{C_+^\delta }\frac{N_{1,n}(\sqrt{\sigma })N_{1,k}(\sqrt{\sigma })^{}}{8\sqrt{\sigma }(\sigma \sigma _k)(\sigma \sigma _n)}e^{2i(\sigma \sigma _n)T_0}𝑑\sigma .\hfill \end{array}$$
(4.55)
The first term is independent of $`T_0`$ (and also of $`\delta >0`$, since the integral converges and the integrand is analytic). In the second term, the real part of the exponent is negative for $`T_0>0`$, so for $`T_0`$ large and positive, the integrand is exponentially small except in a small neighborhood of $`\sigma =0`$. This small neighborhood gives a leading contribution to the integrand that is $`O(T_0^{1/2})`$, and in particular is bounded for large $`T_0>0`$.
Putting these results together, we find that for large $`T_0>0`$, we have
$$_0^{T_0}\gamma (s)𝑑s=(i\overline{M_2}i\mathrm{\Lambda }_2+\mathrm{\Gamma }_2)T_0+O(1).$$
(4.56)
Going back to (4.45), it is clear that choosing
$$_{T_2}D=(i\overline{M_2}i\mathrm{\Lambda }_2+\mathrm{\Gamma }_2)D,$$
(4.57)
will lead to a solution $`h_2(T_0)`$ that is uniformly bounded for all $`T_0>0`$. Also note that the first term in $`\gamma (T_0)`$ contributes a term
$$h_{2,M}i_0^{T_0}\left(M_2(s)\overline{M_2}\right)𝑑s,$$
(4.58)
to the expression for $`h_2(T_0)`$. We write $`h_2(T_0)=h_{2,M}(T_0)+\stackrel{~}{h}_2(T_0)`$.
To find the behavior of $`\gamma (T_0)`$ and its time integral as $`T_0`$ tends to $`\mathrm{}`$, we repeat the above steps, this time deforming the integration contour into the upper half plane to facilitate the steepest descents argument. The path of integration is now $`C_{}^\delta `$ (see Figure 5). The only difference is in the sign of $`\mathrm{\Gamma }_2`$; the correct choice for a bounded solution for all $`T_0`$ is
$$_{T_2}D=(i\overline{M_2}i\mathrm{\Lambda }_2+\mathrm{sgn}(T_0)\mathrm{\Gamma }_2)D.$$
(4.59)
Thus, the method of multiple scales gives the following approximation to the bound state mode amplitude:
$$\begin{array}{c}A_\mathrm{b}(t)=A_{\mathrm{b0}}e^{2i\beta _\mathrm{b}t}e^{i(ϵ\overline{M_1}+ϵ^2\overline{M_2})t}e^{iϵ^2\mathrm{\Lambda }_2t}e^{ϵ^2\mathrm{\Gamma }_2|t|}(1iϵ_0^t(M_1(s)\overline{M_1})ds\hfill \\ \\ \frac{ϵ^2}{2}\left(_0^t(M_1(s)\overline{M_1})ds\right)^2iϵ^2_0^t(M_2(s)\overline{M_2})ds+ϵ^2\stackrel{~}{h}_2(t)+O(ϵ^3)).\hfill \end{array}$$
(4.60)
It is not hard to see that an asymptotically equivalent expression is just
$$A_\mathrm{b}(t)=A_{\mathrm{b0}}e^{2i\beta _\mathrm{b}t}e^{ϵ^2\mathrm{\Gamma }_2|t|}e^{iϵ^2\mathrm{\Lambda }_2t}e^{i_0^tM(s)𝑑s}\left(1+O(ϵ^2)\right).$$
(4.61)
This asymptotic formula is expected to be uniformly valid as $`ϵ`$ tends to zero for all $`|t|<Kϵ^2`$ for any constant $`K`$.
So the behavior of the bound state amplitude under the influence of a periodic perturbation, as predicted by the multiple scale theory, is dominated by two effects, a shift in frequency accompanied by exponential decay. The shift in frequency is an order $`O(ϵ)`$ effect, coming from $`\overline{M}`$. This shift can be traced back to the influence of the perturbation directly on the bound state; there is no coupling to any other modes in this term. The order $`O(ϵ^2)`$ effects include both a further adjustment to the frequency through the quantity $`ϵ^2\mathrm{\Lambda }_2`$, the Lamb shift, and exponential decay through the quantity $`ϵ^2\mathrm{\Gamma }_2`$. Clearly these two numbers are the real and imaginary parts of the same complex frequency. Unlike the leading order phase shift, both of these effects are clearly due to the resonant coupling between the bound state and the continuum that is introduced and mediated by the periodic perturbation. Due to the exponential decay, the lifetime of the bound state is seen to be approximately $`ϵ^2/\mathrm{\Gamma }_2`$, which is quite long for small $`ϵ`$. For this reason, under small perturbations of the potential energy the state is called metastable.
Remark: The validity of this expansion procedure is clearly called into question if any of the resonances $`\sigma _n`$ are very close to zero, in which case the complex frequency $`\mathrm{\Lambda }_2+i\mathrm{\Gamma }_2`$ is potentially large. The breakdown of the expansion in this case indicates the presence of a parametric zero energy resonance. Note, however, that in the odd case the matrix element $`N^{(\mathrm{p})}(t,\lambda )`$ vanishes as $`\lambda `$ tends to zero, and therefore so do the corresponding Fourier coefficients (and in particular they vanish to leading order in $`ϵ`$, that is, $`N_{1,n}(\lambda )`$ vanishes at $`\lambda =0`$ for all $`n`$). This suggests that the expansion (4.61) continues to hold in the odd case as the parameters $`\rho _1`$ and $`\rho _2`$ of the two-soliton potential are varied so as to cause a resonance $`\sigma _n(\rho _1,\rho _2)`$ to pass through zero. In the even case, however, behavior possibly very different from that predicted by the formula (4.61) is expected if a resonance is close to zero. We plan to investigate this phenomenon analytically; however in this paper we will demonstrate the effects of parametric zero energy resonance in both the even and odd cases with numerical simulations. Sudden changes in the behavior of a simple model for atomic ionization as a parameter is smoothly varied, causing the system to pass through a zero energy resonance, have recently been observed and compared with experiment by Costin, Lebowitz and Rokhlenko . $`\mathrm{}`$
### 4.3 Rigorous analysis and infinite time results.
The multiple-scale analysis of the preceding section leads to an asymptotic formula for the decaying bound state amplitude that is valid on time intervals of order $`ϵ^2`$. In this section, we will establish the validity of the asymptotic formula (4.61) in certain circumstances using the results of Kirr and Weinstein . When applicable, these results also yield a detailed description of the solution as $`t\pm \mathrm{}`$.
More precisely, we now study the perturbed periodic system in the form obtained by use of Floquet factorization of the time-periodic unperturbed Hamiltonian $`_0(t)`$:
$$i_tyy=\stackrel{~}{𝒲}(t)y.$$
(4.62)
The self-adjoint operator $`:L^2()L^2()`$ defined in §3 can be thought of as a time-independent Hamiltonian, and the idea is to apply the theory of periodic (or almost periodic) perturbations of autonomous linear Hamiltonian systems as developed in and directly to the problem in this form.
As we did in the multiple scales analysis, we will restrict attention to the special case of periodically perturbed even two-soliton periodic potentials. As we know, in this case the operator $``$ has exactly two $`L^2`$ eigenfunctions, one an even function of $`x`$ and the other an odd function of $`x`$. Since $`L^2`$ is the direct sum of its two subspaces $`L_{(\mathrm{e},\mathrm{o})}^2`$ of even and odd functions, and since $``$ leaves each subspace invariant, we may study the problem (4.62) restricted to one subspace at a time. This reduction results in an unperturbed problem with a single bound state, and to such problems the results described in and can be applied without modification.
On each subspace $`L_{(\mathrm{e},\mathrm{o})}^2()`$, the operator $``$ is explicitly given by:
$$f(x)=\mathrm{\Psi }_\mathrm{b}^{(\mathrm{e},\mathrm{o})}(,0),2\beta _\mathrm{b}f()\mathrm{\Psi }_\mathrm{b}^{(\mathrm{e},\mathrm{o})}(x,0)+_0^{\mathrm{}}\mathrm{\Psi }_\mathrm{d}^{(\mathrm{e},\mathrm{o})}(,0,\lambda ),2\lambda ^2f()\mathrm{\Psi }_\mathrm{d}^{(\mathrm{e},\mathrm{o})}(x,0,\lambda )𝑑\lambda ,$$
(4.63)
where the functions $`\mathrm{\Psi }_\mathrm{b}^{(\mathrm{e},\mathrm{o})}(x,t)`$ and $`\mathrm{\Psi }_\mathrm{d}^{(\mathrm{e},\mathrm{o})}(x,t,\lambda )`$ are defined given the parameters $`\rho _1`$ and $`\rho _2`$ in §2. The hypotheses required in of the even and odd restrictions of the operator $``$ are reproduced here adapted to our application:
* The even and odd restrictions of $``$ are densely defined on subspaces of $`L_{(\mathrm{e},\mathrm{o})}^2()`$ and have self-adjoint extensions to all of $`L_{(\mathrm{e},\mathrm{o})}^2()`$.
* The spectrum of $``$ in each of $`L_{(\mathrm{e},\mathrm{o})}^2()`$ consists of an absolutely continuous part $`\sigma _{\mathrm{cont}}^{(\mathrm{e},\mathrm{o})}()=[0,\mathrm{}]`$ with associated spectral projection $`𝒫_\mathrm{c}^{(\mathrm{e},\mathrm{o})}`$ and a single isolated eigenvalue $`\lambda _0=2\beta _\mathrm{b}`$ with corresponding normalized eigenstate $`\psi _0(x)=\mathrm{\Psi }_\mathrm{b}^{(\mathrm{e},\mathrm{o})}(x,0)`$, so that
$$\psi _0=\lambda _0\psi _0,\psi _0_2=1.$$
(4.64)
* The odd restriction of $``$ satisfies two dispersive local decay estimates. There exist constants $`C_{\mathrm{ns}}`$ and $`C_\mathrm{s}`$ such that
+ The nonsingular local decay estimate
$$^{7/2}e^{it}𝒫_\mathrm{c}^{(\mathrm{o})}f_2C_{\mathrm{ns}}t^{3/2}^{7/2}f_2,$$
(4.65)
holds for all $`fL_{(\mathrm{o})}^2()`$.
+ The singular local decay estimate
$$^{7/2}e^{it}(2\mu 2i\kappa 0)^1𝒫_\mathrm{c}^{(\mathrm{o})}f_2C_\mathrm{s}t^{3/2}^{7/2}f_2,$$
(4.66)
where $`\kappa =\mathrm{sgn}(t)`$, holds uniformly for all $`\mu `$ satisfying $`|\mu |\mu _{\mathrm{min}}>0`$, that is, the constant $`C_\mathrm{s}`$ only depends on $`\mu _{\mathrm{min}}`$.
The local decay hypotheses are established in Appendix B. We remark here that due to a zero energy resonance, the decay estimates that are established in Appendix B for the even case are of the form (4.65) and (4.66) but with decay rate $`t^{1/2}`$ (this is a sharp estimate). Unfortunately, this slower rate of decay precludes the direct application of the results in and . On the other hand, as long as the perturbation does not create a resonance $`\mu `$ that is close to zero, we can expect similar results to hold in the even case over time scales of length $`|t|<K/ϵ^2`$, since there is no obvious difficulty with the multiple scales analysis.
The application of the results of and also requires some hypotheses to be satisfied by the perturbation operator $`\stackrel{~}{𝒲}(t)`$ and its relation to the unperturbed Hamiltonian $``$. The perturbation operator acting on $`L_{(\mathrm{e},\mathrm{o})}^2()`$ takes the form
$$\begin{array}{ccc}\hfill \stackrel{~}{𝒲}(t)f(x)& =& (M(t)\mathrm{\Psi }_\mathrm{b}^{(\mathrm{e},\mathrm{o})}(,0),f()\hfill \\ & & \\ & & +_0^{\mathrm{}}N^{(\mathrm{p})}(t,\eta )\mathrm{\Psi }_\mathrm{d}^{(\mathrm{e},\mathrm{o})}(,0,\eta ),f()d\eta )\mathrm{\Psi }_\mathrm{b}^{(\mathrm{e},\mathrm{o})}(x,0)\hfill \\ & & \\ & & +_0^{\mathrm{}}(N^{(\mathrm{p})}(t,\lambda )^{}\mathrm{\Psi }_\mathrm{b}^{(\mathrm{e},\mathrm{o})}(,0),f()\hfill \\ & & \\ & & +_0^{\mathrm{}}K^{(\mathrm{p})}(t,\lambda ,\eta )\mathrm{\Psi }_\mathrm{d}^{(\mathrm{e},\mathrm{o})}(,0,\eta ),f()d\eta )\mathrm{\Psi }_\mathrm{d}^{(\mathrm{e},\mathrm{o})}(x,0,\lambda )d\lambda .\hfill \end{array}$$
(4.67)
We recall that the periodic “matrix elements” in the above expression are defined in terms of either the odd or even modes by (2.33). This operator, being periodic in $`t`$ with period $`L`$, has a Fourier series expansion
$$\stackrel{~}{𝒲}(t)=\underset{k=\mathrm{}}{\overset{\mathrm{}}{}}e^{2\pi ikt/L}\stackrel{~}{𝒲}_k,$$
(4.68)
where each operator $`\stackrel{~}{𝒲}_k`$ has the same form as (4.67) with the functions $`M(t)`$, $`N^{(\mathrm{p})}(t,\lambda )`$, and $`K^{(\mathrm{p})}(t,\lambda ,\eta )`$ replaced by the corresponding $`k`$th Fourier coefficients, $`M_k`$, $`N_k(\lambda )`$ and $`K_k(\lambda ,\eta )`$ respectively. The hypotheses required in of the perturbation adapted to this context are:
* The operators $`\stackrel{~}{𝒲}_k`$ satisfy $`\stackrel{~}{𝒲}_k=\stackrel{~}{𝒲}_k^{}`$ and
$$\underset{k=\mathrm{}}{\overset{\mathrm{}}{}}\stackrel{~}{𝒲}_k_{(L^2())}<\mathrm{},$$
(4.69)
where $`_{(L^2())}`$ denotes the uniform operator norm in $`L^2()`$. Also,
$$|\stackrel{~}{𝒲}()|\underset{k=\mathrm{}}{\overset{\mathrm{}}{}}\left(^{7/2}\stackrel{~}{𝒲}_k_{(L^2())}+^{7/2}\stackrel{~}{𝒲}_k^{7/2}_{(L^2())}\right)<\mathrm{}.$$
(4.70)
* The following resonance condition holds:
$$\mathrm{\Gamma }\pi \underset{n=n_0}{\overset{\mathrm{}}{}}\stackrel{~}{𝒲}_n\psi _0,\delta (2\sigma _n)\stackrel{~}{𝒲}_n\psi _0>0,$$
(4.71)
where the resonances are defined by $`\sigma _n=(\lambda _0+2\pi n/L)/2=\beta _\mathrm{b}+n\pi /L`$. Here $`n_0`$ is the smallest positive integer for which $`\sigma _{n_0}>0`$. Note that since
$$\stackrel{~}{𝒲}_n\psi _0=\stackrel{~}{𝒲}_n\mathrm{\Psi }_\mathrm{b}^{(\mathrm{e},\mathrm{o})}(x,0)=M_n\mathrm{\Psi }_\mathrm{b}^{(\mathrm{e},\mathrm{o})}(x,0)+_0^{\mathrm{}}N_n(\lambda )^{}\mathrm{\Psi }_\mathrm{d}^{(\mathrm{e},\mathrm{o})}(x,0,\lambda )𝑑\lambda ,$$
(4.72)
and since for $`\sigma >0`$
$$f(),\delta (2\sigma )f()=_0^{\mathrm{}}|(𝒱𝒫_\mathrm{c}^{(\mathrm{e},\mathrm{o})}f)(\lambda )|^2\delta (2\lambda ^22\sigma )𝑑\lambda =\frac{|(𝒱𝒫_\mathrm{c}^{(\mathrm{e},\mathrm{o})}f)(\sqrt{\sigma })|^2}{4\sqrt{\sigma }},$$
(4.73)
The formula for $`\mathrm{\Gamma }`$ can be written as
$$\mathrm{\Gamma }=\frac{\pi }{4}\underset{n=n_0}{\overset{\mathrm{}}{}}\frac{N_n(\sqrt{\sigma _n})}{\sqrt{\sigma _n}}.$$
(4.74)
Note that if $`N_n(\lambda )`$ has an expansion in a small parameter $`ϵ`$, with leading term linear in $`ϵ`$ as was assumed in the multiple scale analysis, then the leading term of the corresponding expansion for $`\mathrm{\Gamma }(ϵ)`$ is exactly $`ϵ^2\mathrm{\Gamma }_2`$, where $`\mathrm{\Gamma }_2`$ is correctly obtained by the multiple scale analysis and is given by (4.54). The constant $`\mathrm{\Gamma }`$ is a decay rate associated with the bound state of the unperturbed system. The statement that the expression (4.71) should be positive for decay to occur as a consequence of resonant coupling to the continuum is attributed to Fermi and is known as the “Fermi Golden Rule”. Again, because the decay constant $`\mathrm{\Gamma }`$ is quadratic in the size of the perturbation, the exponential decay process is very slow for small perturbations. Thus, in the presence of a small perturbation $`W(x,t)`$, the bound state is said to be metastable.
* There are no finite accumulation points of the resonances $`\sigma _n`$, $`nn_0`$. This is satisfied automatically because the Fourier expansion of $`\stackrel{~}{𝒲}(t)`$ is that of a periodic function. The point here is that the results in are more general; for example this hypothesis is satisfied by finite Fourier sums with incommensurate frequencies. Yet further generalizations can be found in .
Verifying the hypothesis (H4) would seem to require more detailed information about the correction to the potential energy $`W(x,t)`$ than we have used thus far. We merely point out at this time that by elementary Cauchy-Schwarz arguments applied to the unitarily equivalent operators $`𝒱\stackrel{~}{𝒲}_n𝒱^{}`$, one finds the estimate
$$\stackrel{~}{𝒲}_n_{(L^2())}2\sqrt{|M_k|^2+2_0^{\mathrm{}}|N_n(\lambda )|^2𝑑\lambda +_0^{\mathrm{}}_0^{\mathrm{}}|K_n(\lambda ,\eta )|^2𝑑\lambda 𝑑\eta }.$$
(4.75)
Assuming these bounds all are finite, which is really a question of the smoothness and decay of “snapshots” of the function $`W(x,t)`$ at fixed $`t`$, we see that the first required bound in (H4) will be satisfied if the Fourier coefficients of the function $`W(x,t)`$ in $`t`$ decay faster than, say, $`1/n`$. This is because the other periodic contributions come from the analytic eigenfunctions, whose Fourier coefficients decay faster than $`1/n^p`$ for any $`p>0`$. Therefore, not much beyond continuity in $`t`$ is required of $`W(x,t)`$, at least for this simpler estimate. More restrictions are certainly required to satisfy the second estimate of (H4).
These hypotheses imply the following results:
###### Proposition 4.1 (Theorem 2.1 of )
Let $``$ and $`\stackrel{~}{𝒲}(t)`$ satisfy the above hypotheses and let an odd function $`y_0(x)`$ be given such that $`x^{7/2}y_0(x)L_{(\mathrm{o})}^2()`$. Let $`y(x,t)`$ be the solution of (4.62) with initial condition $`y(x,0)=y_0(x)`$. Then if $`|\stackrel{~}{𝒲}()|`$ is sufficiently small there exists a constant $`C`$ such that
$$^{7/2}y(,t)_2Ct^{3/2}^{7/2}y_0()_2,$$
(4.76)
holds for all $`t`$.
###### Proposition 4.2 (Theorem 2.2 of )
Assume the same hypotheses of $``$ and $`\stackrel{~}{𝒲}(t)`$. Then if $`|\stackrel{~}{𝒲}()|`$ is sufficiently small, the solution $`y(x,t)`$ of (4.62) corresponding to the odd initial condition $`y_0(x)`$ with $`x^{7/2}y_0(x)L_{(\mathrm{o})}^2()`$ is of the form
$$\begin{array}{ccc}\hfill y(x,t)& =& \left[\mathrm{\Psi }_\mathrm{b}^{(\mathrm{o})}(,0),y_0()e^{2i\beta _\mathrm{b}t}e^{\mathrm{\Gamma }|t|}e^{i\mathrm{\Lambda }t}e^{i_0^tM(s)𝑑s}e^{ir_1(t)}+r_2(t)\right]\mathrm{\Psi }_\mathrm{b}^{(\mathrm{o})}(x,0)\hfill \\ & & \\ & & +\left(e^{it}𝒫_\mathrm{c}^{(\mathrm{o})}y_0()\right)(x,t)+\stackrel{~}{y}(x,t),\hfill \end{array}$$
(4.77)
where
$$\mathrm{\Lambda }\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}\stackrel{~}{𝒲}_n\psi _0,\mathrm{P}.\mathrm{V}.(2\sigma _n)^1𝒫_\mathrm{c}^{(\mathrm{o})}\stackrel{~}{𝒲}_n\psi _0,$$
(4.78)
and where
* The phase correction $`r_1(t)`$ is uniformly bounded and $`O(|\stackrel{~}{𝒲}()|^2)`$ ,
* The bound state amplitude error $`r_2(t)`$ is $`O(|\stackrel{~}{𝒲}()|)`$ uniformly for all $`|t|<K/\mathrm{\Gamma }`$ for all fixed $`K`$ and decays for large time as $`O(t^{3/2})`$,
* The correction $`\stackrel{~}{y}(x,t)`$ is orthogonal to the bound state: $`\psi _0,\stackrel{~}{y}(,t)=0`$ for all $`t`$, and satisfies the dispersive decay estimate $`^{7/2}\stackrel{~}{y}(,t)_2=O(t^{3/2})`$ for large $`t`$.
Remark: Propositions 4.1 and 4.2 would appear to say that all initial conditions decay exponentially and then algebraically. However, a more careful reading shows that it is possible for there to be a transient stage of growth, before the decay ultimately sets in. This is because the error terms, although small when the perturbation is small, are not uniformly small for all initial conditions $`y_0(x)`$ such that $`x^{7/2}y_0(x)`$ ranges over the unit sphere in $`L^2()`$. So, for each fixed perturbation $`W(x,t)`$, no matter how small, it is possible to find an initial condition $`y_0(x,t)`$ that grows before it decays. This is achieved by the following thought experiment. Suppose the periodic perturbation $`W(x,t)`$ is fixed and even in time $`t`$. Now pick any initial condition $`y_0(x)`$ so that $`x^{7/2}y_0(x)L^2()`$. Proposition 4.1 guarantees that after a sufficiently large number $`N`$ of periods, the size of the solution of (4.62) when measured in the weighted $`L^2`$ norm is as small as we please. Note that throughout this process, the solution continues to satisfy $`^{7/2}y(,t)_2<\mathrm{}`$. So now, start again at $`t=0`$ with the new initial condition $`y_0(x)=y(x,t=NL)^{}`$. Since the potential is real and even in time, integration of (4.62) with this new initial condition is, up to complex conjugation, equivalent to integration backwards in time from the time $`t=NL`$ with the initial condition $`y(x,t=NL)`$. So we know that for this very small initial condition, the weighted $`L^2`$ norm must first grow to an order one size at time $`t=NL`$ as the decay process transiently reverses itself, before ultimately giving way to decay over longer times. The existence of such solutions does not violate the statement of Proposition 4.1 or Proposition 4.2 because if one keeps the same initial condition and then makes the perturbation smaller yet again, the connection with the time-reversed problem is lost for this initial condition, and decay occurs sooner. $`\mathrm{}`$
By the same arguments applied in the above discussion of the decay constant $`\mathrm{\Gamma }`$, it follows that there is an alternative formula for $`\mathrm{\Lambda }`$:
$$\mathrm{\Lambda }=\underset{n=\mathrm{}}{\overset{n_01}{}}_0^{\mathrm{}}\frac{|N_n(\sqrt{\sigma })|^2d\sigma }{4\sqrt{\sigma }(\sigma \sigma _n)}\underset{n=n_0}{\overset{\mathrm{}}{}}\mathrm{P}.\mathrm{V}._0^{\mathrm{}}\frac{|N_n(\sqrt{\sigma })|^2d\sigma }{4\sqrt{\sigma }(\sigma \sigma _n)}.$$
(4.79)
Again, if $`N_n(\lambda )`$ has an expansion in a small parameter $`ϵ`$ of the form $`N_n(\lambda )=ϵN_{1,n}(\lambda )+O(ϵ^2)`$ then the leading term of $`\mathrm{\Lambda }`$ is of the form $`ϵ^2\mathrm{\Lambda }_2`$, where $`\mathrm{\Lambda }_2`$ as given by (4.53) was resolved by the multiple scales analysis. This frequency shift associated with the decay of the bound state is the Lamb shift.
From these results, one recovers the true dynamics by setting $`f(x,t)=(𝒫(t)y(,t))(x,t)`$, where $`𝒫(t)`$ is the periodic operator that appeared in the Floquet factorization of the propagator $`𝒰(t)`$ for the periodic unperturbed Hamiltonian $`_0(t)`$. Since
$$(𝒫(t)e^{2i\beta _\mathrm{b}t}\mathrm{\Psi }_\mathrm{b}^{(\mathrm{o})}(,t))(x,t)=(𝒫(t)e^{it}\mathrm{\Psi }_\mathrm{b}^{(\mathrm{o})}(,t))(x,t)=(𝒰(t)\mathrm{\Psi }_\mathrm{b}^{(\mathrm{o})}(,t))(x,t)=\mathrm{\Psi }_\mathrm{b}^{(\mathrm{o})}(x,t),$$
(4.80)
it follows that the time-dependent projection of $`f(x,t)`$ onto the bound state Bloch function $`\mathrm{\Psi }_\mathrm{b}^{(\mathrm{o})}(x,t)`$ is uniformly approximated by
$$B_\mathrm{b}(t)\mathrm{\Psi }_\mathrm{b}^{(\mathrm{o})}(,t),f(,t)\mathrm{\Psi }_\mathrm{b}^{(\mathrm{o})}(,0),f(,0)e^{\mathrm{\Gamma }|t|}e^{i\mathrm{\Lambda }t}e^{i_0^tM(s)𝑑s}.$$
(4.81)
For the system restricted to the odd part of $`L^2()`$, these theorems provide justification for the formal multiple scales analysis carried out above, and more. They globally describe the decay process for all time, where the multiple scales calculation only attempts to capture the dynamics over time scales of length $`\mathrm{\Gamma }^1`$. On the other hand, since the rate of free dispersive decay is not sufficient in the even case to apply this detailed theory, we must settle for the multiple scale expansions.
## 5 Applications in Planar Waveguide Optics
In this section, we present a physical application of the kinds of perturbed time-dependent Schrödinger equations we have been studying in detail. This will provide a concrete family of perturbations $`W(x,t)`$ that we can use in subsequent numerical experiments.
### 5.1 Time-dependent Schrödinger equations in waveguide optics.
For completeness, we present here a brief derivation of the time-dependent Schrödinger equation as it occurs in the paraxial theory of monochromatic waveguide optics. Consider Maxwell’s wave equation for the electric field vector $`\stackrel{}{E}(\stackrel{}{x},t)`$ in a planar ($`\stackrel{}{x}=(y,z)`$) dielectric medium with isotropic, inhomogeneous linear susceptibility $`\chi ^{(1)}(\stackrel{}{x},t)`$
$$\mathrm{\Delta }\stackrel{}{E}\frac{1}{c^2}\stackrel{}{E}_{tt}(\stackrel{}{E})=\frac{1}{c^2}[\chi ^{(1)}(\stackrel{}{x},t)\stackrel{}{E}]_{tt}.$$
(5.1)
Here, the asterisk indicates convolution in time. A Fourier transform (denoted with the operator $``$) in $`t`$ with dual variable $`\omega `$ (the optical frequency) leads to
$$\mathrm{\Delta }\stackrel{}{E}(\stackrel{}{E})+\frac{\omega ^2n^2(y,z,\omega )}{c^2}\stackrel{}{E}=0,$$
(5.2)
where the refractive index $`n`$ is defined by $`n^2(y,z,\omega )1+(\chi ^{(1)})(y,z,\omega )`$. We now assume that the inhomogeneity is weak, so that gradients of $`(\chi ^{(1)})(y,z,\omega )`$ are small. This implies that in the absence of any free charges, the approximate relation $`\stackrel{}{E}0`$ follows from the exact relation for the electric displacement $`\stackrel{}{D}=0`$. Neglecting the divergence term in (5.2), one may then choose any unit vector $`\stackrel{}{e}`$ and set $`(\stackrel{}{E})(y,z,\omega )=\varphi (y,z,\omega )\stackrel{}{e}`$, which gives the Helmholtz or scalar wave equation for $`\varphi `$:
$$\varphi _{zz}+\varphi _{yy}+\frac{\omega ^2n^2(y,z,\omega )}{c^2}\varphi =0.$$
(5.3)
In the design of integrated optical devices, the inhomogeneity in the refractive index is a localized modulation of a “background index” $`n_0(\omega )`$. Choose a fixed length scale $`L_0`$ and nondimensionalize by setting $`z/L_0=\delta ^2Z`$ and $`y/L_0=\delta ^1Y`$, where $`\delta `$ is a dimensionless parameter, and $`Y`$ and $`Z`$ are dimensionless coordinates. Setting $`\varphi (y,z)=f(Y,Z)\mathrm{exp}(i\beta Z/\delta ^2)`$, where $`\beta =L_0\omega n_0(\omega )/c`$ is also dimensionless, one arrives at
$$2i\beta \delta ^2f_Z+f_{ZZ}+\delta ^2f_{YY}+\beta ^2\delta ^4\left[\frac{n^2(YL_0\delta ^1,ZL_0\delta ^2,\omega )}{n_0^2(\omega )}1\right]f=0.$$
(5.4)
With the definition
$$Q(Y,Z;\omega )\frac{1}{2\delta ^2}\left[\frac{n^2(YL_0\delta ^1,ZL_0\delta ^2,\omega )}{n_0^2(\omega )}1\right],$$
(5.5)
we see that the formal limit of $`\delta 0`$ with $`\beta `$ and $`Q(Y,Z;\omega )`$ held fixed yields the paraxial wave equation
$$i\beta f_Z+\frac{1}{2}f_{YY}\beta ^2Q(Y,Z;\omega )f=0.$$
(5.6)
The potential function $`Q(Y,Z;\omega )`$ vanishes as the refractive index approaches its background value $`n_0(\omega )`$, say as $`Y`$ and $`Z`$ go off to infinity (at least in most directions). Given a fixed function $`Q(Y,Z;\omega )`$, we see that the paraxial approximation made here ($`\delta 0`$) is valid if the modulation in the refractive index is weak, slowly varying, and more slowly varying in the $`z`$ direction than in the $`y`$ direction. That is, a fixed function $`Q(Y,Z;\omega )`$ provides an asymptotic description of a family of physical refractive index profiles parametrized by $`\delta 1`$:
$$n^2(y,z,\omega ;\delta )=n_0^2(\omega )2\delta ^2n_0^2(\omega )Q(\delta y/L_0,\delta ^2z/L_0;\omega ).$$
(5.7)
Note that these assumptions about the refractive index justify a posteriori our neglect of the term $`(\stackrel{}{E})`$ in the original wave equation, because in the limit $`\delta 0`$, gradients of $`n^2(y,z,\omega ;\delta )`$ necessarily vanish.
### 5.2 Spectral properties of paraxial waveguides.
In optical waveguide theory, integration (numerical or otherwise) of the linear Schrödinger equation (5.6), also known as the beam propagation method, is one of the main tools for studying the optical properties of “long” planar structures like gradual fiber tapers or channel waveguide junctions, in which backward reflecting waves can be neglected. In this connection, a common problem that arises is the description of the change in behavior of a waveguiding structure as the optical frequency is varied in the neighborhood of some frequency $`\omega _0`$. If the structure $`n^2(y,z,\omega )`$ is one that admits the paraxial approximation, we can use the theory described above as a model. In this case, it is convenient to choose the length scale $`L_0`$ so that at the frequency $`\omega _0`$ we have $`\beta =1`$. With this choice, we think of $`\beta =\beta (\omega )`$ as a function of frequency satisfying $`\beta (\omega _0)=1`$. With the function $`Q(Y,Z;\omega )`$ chosen consistently, the problem becomes one of studying the dependence of solutions of (5.6) on the frequency parameter $`\omega `$ near $`\omega _0`$. With the change of variables $`x=Y\sqrt{\beta (\omega )}`$ and $`t=Z`$, the equation (5.6) takes the form
$$if_t=\left(\frac{1}{2}_x^2+V_0(x,t)\right)f+W(x,t)f,$$
(5.8)
where
$$V_0(x,t)=Q(x,t;\omega _0),$$
(5.9)
and the correction to the potential is given by
$$W(x,t)=\beta (\omega )Q(x/\sqrt{\beta (\omega )},t;\omega )Q(x,t;\omega _0).$$
(5.10)
Setting $`ϵ=\omega /\omega _01`$, we see that $`W(x,t)=W(x,t;ϵ)`$ is uniformly small in $`ϵ`$ if $`Q(Y,Z;\omega )`$ is in the Schwartz space with respect to $`Y`$. We have the expansion
$$W(x,t;ϵ)=ϵ\omega _0\left[\beta ^{}(\omega _0)\left(1\frac{x}{2}_x\right)Q(x,t;\omega _0)+_\omega Q(x,t;\omega _0)\right]+O(ϵ^2),$$
(5.11)
uniformly in $`x`$ and $`t`$. If the frequency range of interest is sufficiently small, then it is often a good approximation to consider the problem to be dispersionless, so that the refractive index $`n(y,z,\omega )`$ is independent of $`\omega `$. In this paper, we will accordingly consider the function $`Q`$ to be independent of $`\omega `$ in which case $`Q(x,t;\omega )=V_0(x,t)`$ for all $`\omega `$ in the range of interest, and we can drop the corresponding term in (5.11).
Suppose now that we choose to study a refractive index profile $`n^2(y,z)`$ that is even in $`y`$ and periodic in $`z`$, such that after choosing a frequency $`\omega _0`$ and nondimensionalizing, the function $`V_0(x,t)`$ is one of the separable potentials described in detail at the end of §2. Over length scales where the paraxial approximation is valid, this periodically modulated channel waveguide will actually have two “breather modes”, approximately described by the bound states $`\mathrm{\Psi }_\mathrm{b}^{(\mathrm{e})}(x,t)`$ and $`\mathrm{\Psi }_\mathrm{b}^{(\mathrm{o})}(x,t)`$. The effect of not being fully in the paraxial limit (that is, $`\delta `$ is small but finite) is that the modes will very slowly attenuate as they propagate forward due to a small coupling to backward propagating fields. This small attenuation occurs at all frequencies near $`\omega =\omega _0`$ in a way that can be quantified . However, the profile $`n^2(y,z)`$ is very special in that at the frequency $`\omega =\omega _0`$ there is no coupling between the bound modes and any forward propagating radiation modes. This additional coupling would indeed be present for “typical” $`z`$-periodic waveguide profiles $`n^2(y,z)`$.
In fact, the theory developed in §4 can be applied to the perturbed problem (5.8) because the unperturbed potential $`V_0(x,t)`$ and the perturbation $`W(x,t;ϵ)`$ are both even functions of $`x`$ that are periodic in $`t`$ with the same period $`L`$. This theory shows that the additional attenuation due to coupling to forward propagating radiation, while completely suppressed at the frequency $`\omega _0`$, reemerges upon detuning the frequency slightly from $`\omega _0`$. Suppose the waveguide is cleaved at $`z=0`$ and is illuminated at this face with a broadband source consisting of many frequencies $`\omega `$. After some distance all of the frequencies will have attenuated somewhat due to backscattering (weak non-paraxiality). However, all frequencies except $`\omega _0`$ will additionally decay by forward propagating radiation damping. The waveguide will therefore preferentially “pass” light of the frequency<sup>8</sup><sup>8</sup>8Actually, the “background” attenuation due to nonparaxiality ($`\delta 0`$) decreases slightly with increasing frequency. When this effect is combined with the frequency-dependent decay calculated from the paraxial approximation, the preferred frequency for which the loss is minimal is detuned slightly upward by an amount that is $`O(\delta ^2)`$. $`\omega _0`$. These effects were observed numerically in .
Note that from the point of view of optical waveguide theory, the periodicity of the index $`n(x,z)`$ in $`z`$ is an important feature, since it gives rise to an attenuated frequency response that is a symmetric function of frequency $`\omega `$ in the neighborhood of $`\omega _0`$. Thus, attenuation occurs whether $`\omega `$ is less than or greater than the frequency $`\omega _0`$ of structural instability. By constrast, channel waveguides, where $`n(y,z)`$ is independent of $`z`$, also exhibit frequency-dependent structural instability at cutoff frequencies where the number of bound states changes. However, in such waveguides the number of bound states (guided modes) is always an increasing function of frequency , which implies that an input beam that matches onto a mode at its cutoff frequency $`\omega _0`$ will attenuate for $`\omega `$ slightly less than $`\omega _0`$ but will remain bound and thus give rise to a significant transmission for $`\omega `$ slightly greater than $`\omega _0`$. Thus, whereas channel waveguides with $`z`$-independent refractive index profiles can behave as “high-pass” components, $`z`$-periodic waveguides that at frequency $`\omega _0`$ are modeled by separable potentials can behave as “band-pass” components.
## 6 Numerical Simulations
Here, we describe some numerical simulations we performed to verify the analytical predictions where we expect them to apply. We also would like to explore the behavior of the perturbed system in parameter regimes where we expect zero-energy resonances (see the remark at the end of §4.2) to prevent the theory from applying in its current form. For concreteness, we considered periodic perturbations of two problems, each associated with a particular two-soliton separable periodic potential. The particular perturbation we selected was exactly the type considered in §5, namely, given a separable two-soliton periodic potential $`V_0(x,t)`$, we numerically integrated the equation
$$i_tf+\frac{1}{2(1+ϵ)}_x^2f(1+ϵ)V_0(x,t)f=0,$$
(6.1)
for several small values of $`ϵ`$. This problem differs from the type to which the theory developed above applies only by a rescaling of $`x`$; in particular, the time scale is unaffected.
Let us give some details about our numerical scheme. We used a Fourier split-step method with a local truncation error of $`O(\mathrm{\Delta }t^3)`$ . The spatial domain $`[x_L,x_R]`$ of $`[80,80]`$ in the “non zero-energy resonance” case and $`[40,40]`$ in the “zero-energy resonance” case (see below for more details about these two cases) was discretized into $`1024`$ points. The scheme splits the Hamiltonian into two parts: $`(t)=_1+_2(t)`$, where
$$_1\frac{1}{2(1+ϵ)}_x^2,_2(t)(1+ϵ)V_0(x,t).$$
(6.2)
Let $`𝒰^ϵ(t,s)`$ denote the propagator associated with (6.1). Let $`𝒰_1^ϵ(ts)`$ and $`𝒰_2^ϵ(t,s)`$ be those associated with $`_1`$ and $`_2(t)`$. Then, the numerical scheme approximates the true integration over a time step of size $`\mathrm{\Delta }t`$ as follows:
$$𝒰^ϵ(t+\mathrm{\Delta }t,t)𝒰_1^ϵ(\mathrm{\Delta }t/4)𝒰_2^ϵ(t+3\mathrm{\Delta }t/4,t+\mathrm{\Delta }t/4)𝒰_1^ϵ(\mathrm{\Delta }t/4),$$
(6.3)
which has an error of order $`\mathrm{\Delta }t^3`$. It is easy to see that, after getting started with a quarter-step, and until finishing with a quarter-step, iterating this approximation to the propagator $`𝒰^ϵ(t,s)`$ over many steps amounts to simply alternating between $`𝒰_1^ϵ`$ and $`𝒰_2^ϵ`$ each acting over a half-step of length $`\mathrm{\Delta }t/2`$.
So, in each half-step, only one of the two parts is integrated. The half-step involving $`_1`$ is carried out in the Fourier transform domain where one multiplies by the explicit exponential of the operator. This step is thus exact in time, so that the only error appears in discretizing the Fourier transform and is smaller than any power of $`\mathrm{\Delta }x`$ if the functions to be differentiated are taken to be arbitrarily smooth. The half-step involving $`_2(t)`$ is done exactly because we have explicit formulas for $`V_0(x,t)`$ and it is possible to find an explicit exponential of $`_2(t)`$. That is, we can write down a formula for the multiplication operator:
$$𝒰_2^ϵ(t,t_0)=\mathrm{exp}\left(i(1+ϵ)_{t_0}^tV_0(x,s)𝑑s\right),$$
(6.4)
and use it in the code. Since the temporal gradients of $`V_0(x,t)`$ can be large in some parts of each period and small in others, we adjusted the time step throughout the period.
We expect the perturbation to generate radiation from the central bound region of the potential, and we need to remove this radiation from the problem as it moves to large $`|x|`$. To take care of this we used a “sponge layer” in which we effectively add a term of the form
$$id\left[\mathrm{exp}\left(\left(\frac{xx_R}{w}\right)^2\right)+\mathrm{exp}\left(\left(\frac{xx_L}{w}\right)^2\right)\right]f,$$
(6.5)
to the right-hand side of (6.1) for a positive damping factor $`d`$ and width $`w`$. These parameters were adjusted heuristically until it was observed, roughly speaking, that no energy was being artificially drawn out of the center and that no energy that was radiated outward was either reflected or transmitted through to the other side of the periodic domain.
We integrated for $`50`$ periods. In all the experiments it was arranged that the fundamental period was $`L=2\pi `$. We initialized the field $`f`$ at $`t=0`$ to be a snapshot of either the odd or the even mode of the unperturbed problem. Then, after integrating, we calculated the projection of the numerical solution onto the exact solution of the unperturbed problem, defining:
$$B_\mathrm{b}(t)\mathrm{\Psi }_\mathrm{b}^{(\mathrm{e},\mathrm{o})}(,t),f(,t).$$
(6.6)
We verified the accuracy of the code by checking that for $`ϵ=0`$ we had $`B_\mathrm{b}(t)1`$ to several digits, even in the presence of the damping in the sponge layer. Note that the function $`B_\mathrm{b}(t)`$ is related to $`A_\mathrm{b}(t)`$ by the simple relation:
$$B_\mathrm{b}(t)=A_\mathrm{b}(t)e^{2i\beta _\mathrm{b}t}.$$
(6.7)
### 6.1 Away from parametric zero-energy resonance.
For the first experiments, we selected $`\rho _1=1/4`$ and $`\rho _2=3/4`$ as the parameters of the function $`V_0(x,t)`$. It is easy to check that the period is $`L=2\pi `$, and that the Floquet exponent of both odd and even bound states may be taken to be $`\beta _\mathrm{b}=\rho _1^2=1/16`$. Therefore the resonances are explicitly given by
$$\sigma _n=\frac{n}{2}\frac{1}{16},$$
(6.8)
none of which are equal to zero. This means that there is no parametric zero-energy resonance, although in the even case there still is a zero-energy resonance corresponding to insufficient dispersive decay. In this case, the formula for the decay constant $`\mathrm{\Gamma }`$ makes sense for both odd and even parity. Furthermore, for odd parity, we have a proof that the asymptotic expansion obtained previously is indeed valid.
In Figure 6, we show plots of $`\mathrm{log}(|B_\mathrm{b}(t)|)`$ for $`ϵ=0.04`$, $`ϵ=0.02`$, and $`ϵ=0.01`$ for an initial condition of odd parity. The numerical results are plotted with solid curves, and superimposed are corresponding graphs of $`\mathrm{\Gamma }|t|`$ calculated from the analytical formula, the analogue of Fermi’s golden rule, and shown with dotted lines.
The main observation here is that the graphs follow the corresponding straight lines, which have slopes that scale like $`ϵ^2`$, as expected. The deviation from the straight lines appears to scale like $`ϵ^2`$ as well, and to decay in time. In Figure 7, we give corresponding plots of the argument of $`B_\mathrm{b}(t)`$ for an initial condition of odd parity.
In these plots, it is easy to see that the phase grows roughly linearly in time, with slope that is $`O(ϵ)`$. This is the contribution to the frequency shift of the term $`\overline{M}`$, which is indeed $`O(ϵ)`$.
Now, we consider an even initial condition, with corresponding projection $`B_\mathrm{b}(t)`$ onto the even mode of the exact solution for $`ϵ=0`$. Figure 8 contains plots of $`\mathrm{log}(|B_\mathrm{b}(t)|)`$ as calculated from the numerical data for $`ϵ=0.04`$, $`ϵ=0.02`$, and $`ϵ=0.01`$ shown in solid curves.
Also plotted are the corresponding decay curves $`\mathrm{\Gamma }|t|`$ shown with dotted lines. Although for even parity there is insufficient dispersive decay for the results of to apply, the decay constant $`\mathrm{\Gamma }`$ (or more precisely as it is obtained in the multiple scale analysis, $`ϵ^2\mathrm{\Gamma }_2`$) is finite because none of the resonances $`\sigma _n`$ are zero, and we see that the multiple scale theory accurately predicts the rate of decay of the bound state even in this case. The plots of the phase of $`B_\mathrm{b}(t)`$ are shown in Figure 9.
Again, one sees that the rate of drift of the phase is $`O(ϵ)`$, as predicted by the multiple scale theory.
The significant new feature apparently contributed by the lack of sufficient dispersive decay for initial conditions of even parity appears to be the quality of the deviations in $`|B_\mathrm{b}(t)|`$ from the “backbone” decay $`e^{\mathrm{\Gamma }|t|}`$. Not only are they larger for fixed $`ϵ`$ than for initial conditions of odd parity, but they have an undulatory character that suggests the possible contribution of subharmonic frequencies to the dynamics. The period of the undulations superimposed on the decay appears to be long compared with $`L`$, the fundamental period of the problem, but also appears to be more or less independent of $`ϵ`$.
### 6.2 At parametric zero-energy resonance.
As a second set of experiments, we considered a potential energy function $`V_0(x,t)`$ obtained from the parameters $`\rho _1=1/\sqrt{2}`$ and $`\rho _2=1`$. In this case, the period is again $`L=2\pi `$, and the Floquet exponents of both modes are $`\beta _\mathrm{b}=\rho _1^2=1/2`$. The corresponding family of resonances is
$$\sigma _n=\frac{n}{2}\frac{1}{2}.$$
(6.9)
One of these values is equal to zero. This condition for parametric zero-energy resonance always goes hand-in-hand with another property of this potential, namely that the Floquet multipliers of both modes are equal to $`1`$. Thus, both odd and even modes are actually periodic functions of $`t`$ with period $`L`$.
At a parametric zero-energy resonance, the dispersive local decay estimates fail to be sufficient to guarantee the applicability of the theory in , also for initial conditions of odd parity. However, in the odd case, the formulas for the decay constant $`\mathrm{\Gamma }`$ and the Lamb shift $`\mathrm{\Lambda }`$ are finite because there is sufficient vanishing in the numerator coming from the missing generalized eigenfunction at $`\sigma =0`$ to cancel and overcome the weaker vanishing of the denominator. Plots of $`\mathrm{log}(|B_\mathrm{b}(t)|)`$ for odd parity corresponding to $`ϵ=0.04`$, $`ϵ=0.02`$, and $`ϵ=0.01`$ are shown in Figure 10 along with dotted lines indicating the analytical prediction of decay.
The prediction of the theory appears to be very accurate indeed. The plots of the phase of $`B_\mathrm{b}(t)`$ are shown in Figure 11.
They show the frequency shift scaling like $`ϵ`$, as we expect from the contribution of the term $`\overline{M}`$. So it appears that for initial conditions of odd parity, there is little if any effect of the parametric zero-energy resonance, although the rate of dispersive decay is smaller here than at more generic parameter values.
Finally, let us examine the behavior of initial conditions of even parity. For such initial conditions and for these parameter values, we have both a simple zero-energy resonance (as one has in the even case for all parameter values) and a parametric zero-energy resonance (as occurs only for very special parameter values). It is easy to see that both $`\mathrm{\Gamma }`$ and $`\mathrm{\Lambda }`$ are infinite in this case, and clearly one cannot expect the multiple scale analysis to be valid. So what can one expect? In Figure 12 we plot $`\mathrm{log}(|B_\mathrm{b}(t)|)`$ for $`ϵ=0.04`$, $`ϵ=0.02`$, and $`ϵ=0.01`$, as before.
This time, rather than superimposing the straight lines $`\mathrm{\Gamma }|t|`$, we might try to compare with a “renormalized” rate of decay given by the formula for $`\mathrm{\Gamma }`$ with the term coming from $`\sigma =0`$ simply dropped. The straight lines calculated from the renormalized version of $`\mathrm{\Gamma }`$ appear dotted on the plots. We still see quite good agreement at the level of a mean drift of $`|B_\mathrm{b}(t)|`$. As in the previous experiment with even parity, we see subharmonic undulations about this mean drift. However, a key point is that whereas previously the period of these undulations appeared to be more or less indpendent of $`ϵ`$, in this case we note that the period appears to scale like $`ϵ^1`$. Thus, there is a “slow” dynamical process involving variations of the amplitude that is completely missed by the multiple scales analysis in its current form. We must expect that whatever rescalings are required to balance the blowing up of $`\mathrm{\Gamma }`$ in the vicinity of a parametric zero-energy resonance will also introduce interesting subharmonic dynamics on the scale of $`T_1=ϵt`$ that will reproduce the effects we are seeing numerically. As a final remark, the phase of $`B_\mathrm{b}(t)`$, as shown in Figure 13
exhibits no particularly different behavior than was seen in any of the other experiments. The frequency adjustment continues to be dominated by the relatively large term $`\overline{M}`$ and is therefore order $`ϵ`$.
## 7 Conclusions
In studying the propagation of waves in time-periodic potentials, considering the problem at hand to be a perturbation of a separable periodic problem is evidently as easy as, and in many cases more convenient than working with periodic perturbations of stationary potential problems. A particular application to the theory of periodically modulated optical waveguides in planar dielectric media allows one to study frequency dependent attenuation properties of certain optical waveguides.
Many of the difficulties described in our paper concern the influence of zero-energy resonances. These are generically not present (that is, for most separable periodic potentials, as for most stationary potentials), but are always present when the potential has sufficient symmetry, as in the evenness considered above. Many problems would therefore vanish upon dropping the symmetry. From one point of view, this introduces the additional complication of having multiple bound states that are essentially coupled to one another by the perturbation. The study of perturbed multimode problems arises naturally in the theory of light propagation in optical fibers having large effective cross-sections. Some of the necessary modifications in the theory described in are described by the same authors in .
Of course another point of view is to keep the symmetry, and hence the possibility of zero-energy resonance, and study the effect of the resonance in more detail. Our numerical experiments suggest that the effects of such a resonance are most dramatic when the expressions for $`\mathrm{\Gamma }`$ and $`\mathrm{\Lambda }`$ blow up, but we also see significant effects, presumably coming simply from the lack of sufficient long time decay of freely dispersing waves, when these quantities are finite. An asymptotic perturbation theory for small $`ϵ`$ should be uniformly valid with respect to parameters like $`\rho _1`$ and $`\rho _2`$, and we plan to investigate zero-energy resonances with such a goal in mind in future work.
## 8 Acknowledgements
P. D. Miller is grateful for the support of the NSF under grant number DMS 9304580 while at the Institute for Advanced Study. A. Soffer is supported in part by a FAS-Rutgers grant and by NSF grant number DMS 9706780. M. I. Weinstein is supported in part by NSF grant number DMS 9500997. Collaboration began while M. I. Weinstein visited the IAS in March 1998 as part of the Program in Geometric Partial Differential Equations organized by Karen Uhlenbeck, whom the authors thank for her support.
## Appendix Appendix A Some theory of separable potentials.
For completeness, we here give a self-contained description of the separable potentials for the linear Schrödinger equation that are connected with the soliton theory of vector nonlinear Schrödinger equations. However, the material is auxiliary and all needed facts are reproduced in the main text. The results here are not new but some arguments are carried out here in more detail.
Let $`q_1(x,t),\mathrm{},q_N(x,t)`$ be given smooth bounded complex functions of real $`x`$ and $`t`$, and let $`A`$ be the vector space of differentiable $`^{N+1}`$-valued functions of $`x`$ and $`t`$. Let $`\lambda `$ be a complex parameter, and consider the two linear operators acting in $`A`$:
$$\begin{array}{ccc}\hfill 𝒳(\lambda ,\stackrel{}{q})& & _x\left[\begin{array}{cc}2i\lambda & \stackrel{}{q}^T\\ \stackrel{}{q}^{}& \mathrm{𝟎}\end{array}\right]\hfill \\ & & \\ & =& _x(2i\lambda 𝐄+𝐔(\stackrel{}{q})),\hfill \end{array}$$
(A.1)
where $`𝐄`$ is a matrix whose elements are given by $`E_{ij}=\delta _{i1}\delta _{j1}`$ and
$$𝐔(\stackrel{}{q})=\left[\begin{array}{cc}0& \stackrel{}{q}^T\\ \stackrel{}{q}^{}& \mathrm{𝟎}\end{array}\right],$$
(A.2)
and
$$\begin{array}{ccc}\hfill 𝒯(\lambda ,\stackrel{}{q})& & _t\left[\begin{array}{cc}2i\lambda ^2+i\stackrel{}{q}^T\stackrel{}{q}^{}/2& \lambda \stackrel{}{q}^T+i_x\stackrel{}{q}^T/2\\ \lambda \stackrel{}{q}^{}+i_x\stackrel{}{q}^{}/2& i\stackrel{}{q}^{}\stackrel{}{q}^T/2\end{array}\right]\hfill \\ & & \\ & =& _t\left(2i\lambda ^2𝐄+\lambda 𝐔(\stackrel{}{q})+\frac{i}{2}𝐕(\stackrel{}{q})\right),\hfill \end{array}$$
(A.3)
where
$$𝐕(\stackrel{}{q})=\left[\begin{array}{cc}\stackrel{}{q}^T\stackrel{}{q}^{}& _x\stackrel{}{q}^T\\ _x\stackrel{}{q}^{}& \stackrel{}{q}^{}\stackrel{}{q}^T\end{array}\right].$$
(A.4)
Here $`\stackrel{}{q}`$ denotes the column vector of the functions $`q_k(x,t)`$ and $`\mathrm{𝟎}`$ denotes the $`N\times N`$ zero matrix. Along with these two operators, we consider their nullspaces, $`K_𝒳(\lambda ,\stackrel{}{q})A`$ and $`K_𝒯(\lambda ,\stackrel{}{q})A`$. For generic $`\lambda `$, these subspaces are $`N+1`$-dimensional, and if restricted to generic fixed $`x`$, $`t`$, and $`\lambda `$ span $`^{N+1}`$.
If the functions $`q_k(x,t)`$ are chosen just right, then the subspaces $`K_𝒳(\lambda ,\stackrel{}{q})`$ and $`K_𝒯(\lambda ,\stackrel{}{q})`$ may coincide for all complex $`\lambda `$: $`K_𝒳=K_𝒯K`$. If this is the case, then the common nullspace will certainly be contained in the nullspace of the commutator: $`KK_{[𝒳,𝒯]}`$. As is easily checked, the commutator $`[𝒳,𝒯]`$ is not a differential operator, but is merely a matrix multiplication operator, with entries depending on $`x`$ and $`t`$ through the $`q_k(x,t)`$. Since the kernel of the commutator contains a subspace $`K`$ of dimension $`N+1`$ for most $`\lambda `$, $`x`$, and $`t`$, this implies that the operators $`𝒳`$ and $`𝒯`$ commute. It is easily checked that the compatibility condition $`[𝒳,𝒯]=\mathrm{𝟎}`$ is equivalent to the vector nonlinear Schrödinger equation:
$$i_t\stackrel{}{q}+\frac{1}{2}_x^2\stackrel{}{q}+(\stackrel{}{q}^T\stackrel{}{q}^{})\stackrel{}{q}=0.$$
(A.5)
It is therefore necessary that (A.5) be satisfied by the functions $`q_k(x,t)`$ if we are to have a basis of simultaneous nullvectors in the common nullspace $`K`$. When they exist, we can collect all these linearly independent column vectors into a square matrix $`𝐅(x,t,\lambda )`$. These ideas admit a natural geometric interpretation in the trivial frame bundle $`E^2`$ with fiber $`GL(N+1,)`$. Here, $`𝒳`$ and $`𝒯`$ are covariant derivative operators for $`E`$ in the $`x`$ and $`t`$ directions, and the condition $`[𝒳,𝒯]=\mathrm{𝟎}`$ means that the curvature of the affine connection specified by $`𝒳`$ and $`𝒯`$ is zero. This implies the existence of parallel global sections $`𝐅(x,t,\lambda )`$ of the bundle $`E`$.
Finding a global section $`𝐅(x,t,\lambda )`$ of $`E`$ given $`\stackrel{}{q}(x,t)`$, (that is, a matrix of simultaneous solution vectors) is not always easy and for this reason, we will adopt a different point of view below. However, it is clear from (A.1) and (A.3) that, given bounded functions $`q_k(x,t)`$ satisfying (A.5) it is possible to develop an asymptotic expansion for $`𝐅(x,t,\lambda )`$ in the limit $`\lambda \mathrm{}`$. The expansion may be sought in the form:
$$𝐅(x,t,\lambda )=\left(c𝕀_{N+1}+\lambda ^1𝐅^{(1)}(x,t)+\lambda ^2𝐅^{(2)}(x,t)+\mathrm{}\right)\left[\begin{array}{cc}e^{2i(\lambda x+\lambda ^2t)}& \stackrel{}{0}^T\\ & \\ \stackrel{}{0}& 𝕀_N\end{array}\right].$$
(A.6)
Here, $`𝕀_D`$ denotes the $`D\times D`$ identity matrix, and $`c`$ is a complex constant. The coefficient matrices $`𝐅^{(p)}(x,t)`$ are determined recursively in terms of $`q_1(x,t),\mathrm{},q_N(x,t)`$ and the constant $`c`$ by collecting powers of $`\lambda `$ in the compatible equations $`𝒳𝐅=𝒯𝐅=0`$. There is some ambiguity in this expansion procedure entering as integration constants at each order. However, it is easy to see that
$$F_{1,k+1}^{(1)}(x,t)=\frac{c}{2i}q_k(x,t),k=1,\mathrm{},N.$$
(A.7)
regardless of the values of the integration constants.
The implications of this compatible structure for linear Schrödinger equations that are of interest to us in this paper are easily stated.
###### Proposition A.1
Suppose that (A.5) is satisfied, and let $`\stackrel{}{v}(x,t,\lambda )K`$ be any simultaneous nullvector of $`𝒳(\lambda ,\stackrel{}{q})`$ and $`𝒯(\lambda ,\stackrel{}{q})`$. Let $`𝒫:A(x,t)`$ be the operator of projection onto the first component. Define the self-consistent potential
$$V_0(x,t)\stackrel{}{q}(x,t)^T\stackrel{}{q}(x,t)^{},$$
(A.8)
and set $`f(x,t,\lambda )=𝒫\stackrel{}{v}(x,t,\lambda )`$. Then it follows that
$$i_tf+\frac{1}{2}_x^2fV_0(x,t)f=0.$$
(A.9)
So, for each complex $`\lambda `$, the function $`f(x,t,\lambda )`$ is a solution of the linear, time-dependent Schrödinger equation with potential (A.8). Solutions corresponding to different values of $`\lambda `$ are linearly independent. Given functions $`q_k(x,t)`$ satisfying the nonlinear system (A.5), one can look to the common nullspace $`K`$ of the linear operators $`𝒳(\lambda ,\stackrel{}{q})`$ and $`𝒯(\lambda ,\stackrel{}{q})`$ as a source of many solutions of the linear equation (A.9).
Remark: Let us try to put these facts in a larger context, and incidentally give the proof of Proposition A.9. It is part of the lore of integrable systems theory that linearized evolution equations connected with integrable systems are solvable in terms of “squared eigenfunctions” coming from the auxiliary linear problems making up the Lax pair for the integrable system. The integrable system (A.5) is the compatibility condition for the equations $`𝒳𝐅=\mathrm{𝟎}`$ and $`𝒯𝐅=\mathrm{𝟎}`$. By a change of variables (gauge transformation) $`𝐅=𝐆\mathrm{exp}(i\lambda xi\lambda ^2t)`$ the two equations take the more familiar form of the Lax pair for (A.5) :
$$_x𝐆=\mathrm{𝐀𝐆},_t𝐆=\mathrm{𝐁𝐆},$$
(A.10)
where
$$𝐀=\left[\begin{array}{cc}i\lambda & \stackrel{}{q}^T\\ \stackrel{}{q}^{}& i\lambda 𝕀\end{array}\right],𝐁=\left[\begin{array}{cc}i\lambda ^2+i\stackrel{}{q}^T\stackrel{}{q}^{}/2& \lambda \stackrel{}{q}^T+i_x\stackrel{}{q}^T/2\\ \lambda \stackrel{}{q}^{}+i_x\stackrel{}{q}^{}/2& i\lambda ^2𝕀i\stackrel{}{q}^{}\stackrel{}{q}^T/2\end{array}\right].$$
(A.11)
If $`𝐆_\alpha `$ and $`𝐆_\beta `$ are any two simultaneous matrix solutions of the Lax pair (A.10), and if $`𝐂`$ is any constant (that is, $`x`$ and $`t`$ independent) matrix, then by setting $`𝐐=𝐆_\alpha \mathrm{𝐂𝐆}_\beta ^1`$, one easily obtains the equations
$$_x𝐐=[𝐀,𝐐],_t𝐐=[𝐁,𝐐].$$
(A.12)
Equations of this form are called Lax equations, and the elements of $`𝐐`$ are the “squared eigenfunctions”. The terminology becomes accurate in the scalar case $`N=1`$ when $`𝐀`$ and $`𝐁`$ are in the Lie algebra $`sl(2)`$. In this case the solutions $`𝐆`$ of the Lax pair can be normalized to be in the Lie group $`SL(2)`$ and therefore have determinant one. Then, because $`𝐆_\beta `$ is $`2\times 2`$ with determinant one, the elements of $`𝐐`$ are seen to be bona fide quadratic forms in the solutions of the Lax pair (A.10). The emphasis in the literature on the $`sl(2)`$-specific terminology of “squared eigenfunctions” for the forms that satisfy the Lax equations (A.12) no doubt bears witness to the fact that so many of the famous integrable equations (e.g. Korteweg-de Vries, scalar nonlinear Schrödinger, sine-Gordon) are associated with $`sl(2)`$ representations.
If one introduces the splitting of a matrix into blocks: $`𝐌=𝐌^\mathrm{D}+𝐌^{\mathrm{OD}}`$ where $`𝐌^\mathrm{D}`$ consists of the $`1\times 1`$ and $`N\times N`$ diagonal blocks of $`𝐌`$, and $`𝐌^{\mathrm{OD}}`$ consists of the $`1\times N`$ and $`N\times 1`$ off-diagonal blocks of $`𝐌`$, and if one introduces $`𝐀_0=𝐀|_{\lambda =0}`$ and $`𝐁_0=𝐁|_{\lambda =0}`$, then it is an exercise to check that the equations (A.12) imply
$$\left[\begin{array}{cc}i& \stackrel{}{0}^T\\ \stackrel{}{0}& i𝕀\end{array}\right]_t𝐐^{\mathrm{OD}}+\frac{1}{2}_x^2𝐐^{\mathrm{OD}}\left[\begin{array}{cc}i& \stackrel{}{0}^T\\ \stackrel{}{0}& i𝕀\end{array}\right][𝐁_0^\mathrm{D},𝐐^{\mathrm{OD}}]\frac{1}{2}[𝐀_0^{\mathrm{OD}},[𝐀_0^{\mathrm{OD}},𝐐^{\mathrm{OD}}]]=\mathrm{𝟎}.$$
(A.13)
If one writes
$$𝐐^{\mathrm{OD}}=\left[\begin{array}{cc}0& \stackrel{}{g}^T\\ \stackrel{}{h}& \mathrm{𝟎}\end{array}\right],$$
(A.14)
then one finds
$$\begin{array}{ccccccccccc}\hfill i_t\stackrel{}{g}^T& +& \frac{1}{2}_x^2\stackrel{}{g}^T& +& \stackrel{}{q}^T\stackrel{}{q}^{}\stackrel{}{g}^T& +& \stackrel{}{q}^T\stackrel{}{h}\stackrel{}{q}^T& +& \stackrel{}{g}^T\stackrel{}{q}^{}\stackrel{}{q}^T& =& \stackrel{}{0}\hfill \\ & & & & & & & & & & \\ \hfill i_t\stackrel{}{h}& +& \frac{1}{2}_x^2\stackrel{}{h}& +& \stackrel{}{q}^{}\stackrel{}{q}^T\stackrel{}{h}& +& \stackrel{}{q}^{}\stackrel{}{g}^T\stackrel{}{q}^{}& +& \stackrel{}{h}\stackrel{}{q}^T\stackrel{}{q}^{}& =& \stackrel{}{0}.\hfill \end{array}$$
(A.15)
These linear equations for $`\stackrel{}{g}`$ and $`\stackrel{}{h}`$ are consistent with the constraint $`\stackrel{}{h}=\stackrel{}{g}^{}`$ at which point it becomes clear that $`\stackrel{}{g}(x,t)`$ satisfies the linearization of the vector nonlinear Schrödinger equation (A.5) about a solution $`\stackrel{}{q}(x,t)`$.
Consider now a particular solution $`\stackrel{}{q}(x,t)`$ of (A.5) and by adjoining a new trivial component $`q_{N+1}(x,t)0`$, view it as a solution $`\stackrel{}{q}^{}(x,t)`$ of (A.5) in the $`N+1`$ component case. From (A.15) it is easily seen that the corresponding components $`g_{N+1}(x,t)`$ and $`h_{N+1}(x,t)`$ satisfy
$$\begin{array}{ccccccc}\hfill i_tg_{N+1}& +& \frac{1}{2}_x^2g_{N+1}& +& (\stackrel{}{q}^T\stackrel{}{q}^{})g_{N+1}& =& 0,\hfill \\ & & & & & & \\ \hfill i_th_{N+1}& +& \frac{1}{2}_x^2h_{N+1}& +& (\stackrel{}{q}^T\stackrel{}{q}^{})h_{N+1}& =& 0,\hfill \end{array}$$
(A.16)
where we have used the fact that $`\stackrel{}{q}^T\stackrel{}{q}^{}=\stackrel{}{q}^T\stackrel{}{q}^{}`$. Now considering the Lax pair (A.10) for the primed potentials, it is easy to see that there exists a nontrivial column vector solution of both equations of the form $`(\stackrel{}{0}_{N+1}^T,\mathrm{exp}(i\lambda x+i\lambda ^2t))^T`$, and that further column vector solutions can then be chosen to have a vanishing last component. Taking the last column of the matrix solution $`𝐆_\beta `$ to be this particular solution, and the first $`N+1`$ columns all to have zeros in the final component, we see that $`𝐆_\beta `$ may be inverted in two independent blocks, and therefore a solution of the linearized equation is given by
$$g_{N+1}(x,t)=Q_{1,N+2}(x,t,\lambda )=\mathrm{exp}(i\lambda xi\lambda ^2t)\underset{k=1}{\overset{N+2}{}}C_{k,N+2}(\lambda )G_{\alpha ,1,k}(x,t,\lambda ).$$
(A.17)
Since the matrix $`𝐂`$ is arbitrary, we may view the sum above as the first component of an arbitrary column vector solution of the Lax pair (A.10) with the primed potentials $`\stackrel{}{q}^{}(x,t)`$. Moreover, since $`q_{N+1}^{}(x,t)0`$, the first component of a solution of the primed Lax pair is also the first component of a solution of the unprimed Lax pair for the fully nontrivial potential $`\stackrel{}{q}(x,t)`$. Reversing the gauge transformation between solutions $`𝐆`$ of the unprimed Lax pair (A.10) and solutions $`𝐅`$ of $`𝒳𝐅=𝒯𝐅=\mathrm{𝟎}`$ then establishes the connection with Proposition A.9.
So, the procedure we are using for solving the time-dependent linear Schrödinger equation is exactly the “squared eigenfunction” linearization of a certain $`N+1`$ component nonlinear Schrödinger equation about a particular solution having $`q_{N+1}^{}(x,t)0`$. The “squared eigenfunctions” solving the linearized problem appear to be linear in this special case because for $`q_{N+1}^{}0`$ the primed Lax pair becomes partly trivial, and the contribution of this trivial part to the matrix $`𝐐`$ is completely explicit (the exponential function that we remove with a gauge transformation). $`\mathrm{}`$
We now return to the construction of self-consistent potentials and the corresponding solutions of (A.9). The nonlinear equation (A.5) is an integrable system by virtue of its representation as the compatibility condition of two linear problems. So there are many well-known ways to find functions $`q_k(x,t)`$ for which the corresponding linear Schrödinger equation can be solved. But as we are interested as much in the common nullspace of $`𝒳`$ and $`𝒯`$ as in the functions $`q_k(x,t)`$, we will now describe an effective approach to finding both at the same time. In this approach, the object of fundamental importance is the common nullspace $`K`$ itself. We construct it first, with the functions $`q_k(x,t)`$ being chosen after the fact precisely so that for any basis matrix $`𝐅`$ of $`K`$, we will have $`𝒳(\lambda ,\stackrel{}{q})𝐅=𝒯(\lambda ,\stackrel{}{q})𝐅=\mathrm{𝟎}`$.
What we know about $`K`$ is that whenever it exists by virtue of the compatibility condition, the assumption that the functions $`q_k(x,t)`$ are bounded (this will be justified below) leads to expansions for large $`\lambda `$ of a basis for $`K`$ of the form (A.6). These expansions are generally only asymptotic; there is no guarantee that there exists a choice of the integration constants such that the expansion (A.6) converges for any $`\lambda `$ at all. However, we now suppose that there exist solutions $`q_k(x,t)`$ of the nonlinear system (A.5) for which an expansion (A.6) not only converges in some deleted neighborhood of $`\lambda =\mathrm{}`$, but actually truncates. For such solutions $`q_k(x,t)`$, if they exist, a basis of the subspace $`K`$ is given exactly by an expression of the form
$$𝐅(x,t,\lambda )=\left(c\lambda ^M𝕀_{N+1}+\underset{p=0}{\overset{M1}{}}\lambda ^p𝐅^{(p)}(x,t)\right)\left[\begin{array}{cc}e^{2i(\lambda x+\lambda ^2t)}& \stackrel{}{0}^T\\ & \\ \stackrel{}{0}& 𝕀_N\end{array}\right],$$
(A.18)
for some positive integer $`M`$, where $`c`$ is a complex constant. We have multiplied by an explicit factor of $`\lambda ^M`$ to bring the sum into polynomial form.
Since we are not considering the functions $`q_k(x,t)`$ to be known, we do not have the option of solving for the coefficient matrices $`𝐅^{(p)}(x,t)`$ by substitution into the equations $`𝒳𝐅=𝒯𝐅=0`$. We therefore must consider them to be arbitrary functions of $`x`$ and $`t`$ until we know otherwise. Without any constraints on the coefficients, we see that the differentiable matrix functions of $`x`$, $`t`$, and $`\lambda `$ of the form (A.18), for given integer values of $`M`$ and $`N`$, form a vector space $`\mathrm{\Lambda }_{N,M}`$ over the complex numbers.
The space $`\mathrm{\Lambda }_{N,M}`$ is very large. If our claim — that appropriate solutions $`q_k(x,t)`$ of the nonlinear system (A.5) exist — is not vacuous, then $`\mathrm{\Lambda }_{N,M}`$ should contain many proper subspaces that may be identified with the common nullspace $`K`$ of $`𝒳(\lambda ,\stackrel{}{q})`$ and $`𝒯(\lambda ,\stackrel{}{q})`$ for some $`\stackrel{}{q}`$. If $`𝐅(x,t,\lambda )`$ is of the form (A.18) and is a basis matrix of one of these subspaces, then it must be determined modulo the constant $`c`$. This means that each such subspace of $`\mathrm{\Lambda }_{N,M}`$ should ultimately be isomorphic to $``$, with the isomorphism being established via the constant $`c`$.
We prepare to isolate the appropriate subspaces of $`\mathrm{\Lambda }_{N,M}`$ by defining a set of discrete data. Let $`𝒟`$ denote an $`M`$-tuple of pairs $`(\lambda _k,\stackrel{}{g}^{(k)})`$ where the $`\lambda _k`$ are distinct numbers in the complex upper half-plane and where the $`\stackrel{}{g}^{(k)}`$ are vectors in $`^N`$. From each vector $`\stackrel{}{g}^{(k)}`$, we build $`N+1`$ vectors in $`^{N+1}`$ by setting
$$\stackrel{}{a}^{(k)}=(1,g_1^{(k)},g_2^{(k)},\mathrm{},g_N^{(k)})^T^{N+1},$$
(A.19)
and for $`j=1,\mathrm{},N`$,
$$\stackrel{}{b}^{(k,j)}=(g_j^{(k)},\stackrel{}{e}_j^T)^T^{N+1},$$
(A.20)
where $`\stackrel{}{e}_j`$ are the usual unit vectors in $`^N`$.
###### Definition A.1
$`\mathrm{\Lambda }_{N,M}^𝒟`$ is the subspace of $`\mathrm{\Lambda }_{N,M}`$ whose elements $`𝐅(x,t,\lambda )`$ satisfy the homogeneous linear conditions
$$𝐅(x,t,\lambda _k)\stackrel{}{a}^{(k)}=\stackrel{}{0},$$
(A.21)
for $`k=1,\mathrm{},M`$ and
$$𝐅(x,t,\lambda _k^{})\stackrel{}{b}^{(k,j)}=\stackrel{}{0},$$
(A.22)
for $`k=1,\mathrm{},M`$ and $`j=1,\mathrm{},N`$.
It is not hard use dimension counting arguments to prove the following:
###### Proposition A.2
Let the discrete data $`𝒟`$ be given. The set of solutions of (A.21) and (A.22) forms a one-dimensional linear subspace of $`\mathrm{\Lambda }_{N,M}`$. The general solution of (A.21) and (A.22) is given by the one-parameter family of matrices (A.18), indexed by the complex parameter $`c`$. Thus, $`\mathrm{\Lambda }_{N,M}^𝒟`$, with the isomorphism being established via the complex constant $`c`$. In particular, if $`c`$ is given, then the coefficient functions $`𝐅^{(p)}(x,t)`$ are uniquely determined as functions of $`x`$ and $`t`$, and if $`c`$ is chosen to be zero, then $`𝐅(x,t,\lambda )`$ is the zero matrix.
This proposition allows us to index the elements of $`\mathrm{\Lambda }_{N,M}^𝒟`$ by the constant $`c`$ which is now a genuine coordinate for the one-dimensional subspace $`\mathrm{\Lambda }_{N,M}^𝒟`$. We indicate the dependence by writing $`𝐅_{𝒟,c}(x,t,\lambda )`$ for the matrices in this subspace.
This proposition is true even if homogeneous constraints less structured than (A.21) and (A.22) are imposed. In order for the dimension count to come out right it is sufficient to choose $`M(N+1)`$ arbitrary complex numbers $`\lambda _k`$ along with corresponding constant vectors $`\stackrel{}{c}^{(k)}^{N+1}`$ (the numbers $`\lambda _k`$ need not all be distinct, as long as the vectors $`\stackrel{}{c}`$ belonging to each $`\lambda _k`$ are linearly independent) and to impose $`𝐅_{𝒟,c}(x,t,\lambda _k)\stackrel{}{c}^{(k)}=\stackrel{}{0}`$ for all $`k=1,\mathrm{},M(N+1)`$. The additional structure in the constraints (A.21) and (A.22) is needed for the following.
###### Proposition A.3
Let discrete data $`𝒟`$ be given, and let $`𝐅_{𝒟,c}(x,t,\lambda )\mathrm{\Lambda }_{N,M}^𝒟`$. Then
$$\frac{1}{c}F_{k+1,1}^{(M1)}(x,t)=\left(\frac{1}{c}F_{1,k+1}^{(M1)}(x,t)\right)^{},k=1,\mathrm{},N.$$
(A.23)
We will have use for this symmetry property below. Its proof is simple.
Proof of Proposition A.23: It is sufficient to consider the case of $`c=1`$, since the coefficient matrices simply scale with $`c`$. It will be convenient to introduce the block form of the coefficient matrices:
$$𝐅^{(p)}(x,t)=\left[\begin{array}{cc}a^{(p)}& \stackrel{}{b}^{(p)T}\\ \stackrel{}{c}^{(p)}& 𝐃^{(p)}\end{array}\right],$$
(A.24)
where $`a^{(p)}(x,t)`$ is a scalar, $`\stackrel{}{b}^{(p)}(x,t)`$ and $`\stackrel{}{c}^{(p)}(x,t)`$ are $`N`$-component vectors, and $`𝐃^{(p)}(x,t)`$ is an $`N\times N`$ matrix. We will prove the stronger result that for all $`p=0,\mathrm{},M1`$,
$$\stackrel{}{c}^{(p)}(x,t)=\stackrel{}{b}^{(p)}(x,t)^{}.$$
(A.25)
In this form, the equations (A.21) and (A.22) take the form of the system:
$$\begin{array}{ccc}\hfill \lambda _k^M+\underset{p=0}{\overset{M1}{}}\lambda _k^p\left(\stackrel{}{b}^{(p)T}\stackrel{}{g}^{(k)}e^{2i(\lambda _kx+\lambda _k^2t)}a^{(p)}\right)& =& 0,\hfill \\ & & \\ \hfill \lambda _k^M\stackrel{}{g}^{(k)}+\underset{p=0}{\overset{M1}{}}\lambda _k^p\left(𝐃^{(p)}\stackrel{}{g}^{(k)}e^{2i(\lambda _kx+\lambda _k^2t)}\stackrel{}{c}^{(p)}\right)& =& \stackrel{}{0},\hfill \\ & & \\ \hfill \lambda _k^M\stackrel{}{g}^{(k)}+\underset{p=0}{\overset{M1}{}}\lambda _k^p\left(a^{(p)}\stackrel{}{g}^{(k)}+e^{2i(\lambda _k^{}x+\lambda _k^2t)}\stackrel{}{b}^{(p)}\right)& =& \stackrel{}{0},\hfill \\ & & \\ \hfill \lambda _k^M𝕀_N+\underset{p=0}{\overset{M1}{}}\lambda _k^p\left(e^{2i(\lambda _k^{}x+\lambda _k^2t)}\stackrel{}{c}^{(p)}\stackrel{}{g}^{(k)T}+𝐃^{(p)}\right)& =& \mathrm{𝟎},\hfill \end{array}$$
(A.26)
where $`k=1,\mathrm{},M`$. From the first and fourth equations, we can eliminate $`a^{(p)}(x,t)`$ and $`𝐃^{(p)}(x,t)`$, $`p=0,\mathrm{},M1`$ in favor of the $`\stackrel{}{b}^{(p)}(x,t)`$ and $`\stackrel{}{c}^{(p)}(x,t)`$. This involves introducing the elements of the inverse $`𝐖`$ of the Vandermonde matrix $`𝐕`$ having elements $`V_{jk}\lambda _j^{k1}`$, but it leads to two decoupled linear systems, one for the $`\stackrel{}{b}^{(p)}(x,t)`$ and the other for the $`\stackrel{}{c}^{(p)}(x,t)`$. These systems are:
$$\begin{array}{ccc}\hfill \underset{r=1}{\overset{M}{}}H_{kr}\stackrel{}{c}^{(r1)}& =& h_k,\hfill \\ & & \\ \hfill \underset{r=1}{\overset{M}{}}H_{kr}^{}\stackrel{}{b}^{(r1)}& =& h_k^{},\hfill \end{array}$$
(A.27)
where
$$H_{kr}V_{kr}e^{2i(\lambda _kx+\lambda _k^2t)}+\underset{s=1}{\overset{M}{}}V_{ks}\underset{j=1}{\overset{M}{}}W_{sj}^{}V_{jr}^{}e^{2i(\lambda _j^{}x+\lambda _j^2t)}\stackrel{}{g}^{(j)T}\stackrel{}{g}^{(k)},$$
(A.28)
and
$$h_k\lambda _k^M\stackrel{}{g}^{(k)}\underset{s=1}{\overset{M}{}}\underset{j=1}{\overset{M}{}}V_{ks}W_{sj}^{}\lambda _j^M\stackrel{}{g}^{(k)}.$$
(A.29)
It is then clear that $`\stackrel{}{c}^{(p)}=\stackrel{}{b}^{(p)}`$ for all $`p`$. $`\mathrm{}`$
So, the emphasis has changed with respect to these matrices and their relation to the functions $`q_k(x,t)`$. Rather than determining the coefficient matrices $`𝐅^{(p)}(x,t)`$ from a given set of functions $`q_k(x,t)`$ solving (A.5) by an asymptotic expansion procedure, we are determining them from the discrete data $`𝒟`$ and a choice of the constant $`c`$. If there is to be any consistency, then we must still have relations between the coefficient matrices $`𝐅^{(p)}(x,t)`$ of $`𝐅_{𝒟,c}`$ and the functions $`q_k(x,t)`$; in particular, we can rewrite (A.7):
$$q_k(x,t)\frac{2i}{c}F_{1,k+1}^{(M1)}(x,t)$$
(A.30)
and use it as a definition of some functions $`q_k(x,t)`$ in terms of the discrete data $`𝒟`$ and the constant $`c`$. Note that as long as $`c0`$, then this definition is actually independent of $`c`$ because $`𝐅_{𝒟,c}`$ is directly proportional to $`c`$. The fact that (A.30) is sensible as a definition of the $`q_k(x,t)`$ is shown by:
###### Proposition A.4
Let the discrete data $`𝒟`$ be given and let the constant $`c`$ be nonzero, and let the functions $`q_1(x,t),\mathrm{},q_N(x,t)`$ be defined (in terms of $`𝒟`$ alone) via (A.30). This determines the operators $`𝒳(\lambda ,\stackrel{}{q})`$ and $`𝒯(\lambda ,\stackrel{}{q})`$. Then for any $`𝐅_{𝒟,c}(x,t,\lambda )\mathrm{\Lambda }_{N,M}^𝒟`$,
$$𝒳(\lambda ,\stackrel{}{q})𝐅_{𝒟,c}(\lambda ,x,t)=𝒯(\lambda ,\stackrel{}{q})𝐅_{𝒟,c}(\lambda ,x,t)=\mathrm{𝟎}.$$
(A.31)
For these $`q_k(x,t)`$ the columns of $`𝐅_{𝒟,c}(x,t,\lambda )`$ are generically linearly independent and therefore form a basis of the common nullspace $`K`$ for almost all $`\lambda `$.
Recall that the commutator $`[𝒳,𝒯]`$ is, for fixed $`x`$ and $`t`$, a matrix multiplication operator in $`^{N+1}`$. Thus, the existence of the common nullspace $`K`$ of $`𝒳(\lambda ,\stackrel{}{q})`$ and $`𝒯(\lambda ,\stackrel{}{q})`$ of generic dimension $`N+1`$ for these functions $`q_k(x,t)`$ implies the vanishing of the commutator and the compatibility of the two linear problems. Therefore we have
###### Corollary A.1
The functions $`q_k(x,t)`$ constructed from any set of discrete data $`𝒟`$ satisfy the vector nonlinear Schrödinger equation (A.5).
A time-dependent self-consistent potential function $`V_0(x,t)`$ generated from the functions $`q_k(x,t)`$ connected with a set of discrete data $`𝒟`$ according to (A.8) will be called a separable potential .
Proof of Proposition A.4: Let $`𝐅_{𝒟,c}(x,t,\lambda )\mathrm{\Lambda }_{N,M}^𝒟`$. The proof begins with the simple observation that, as a consequence of the vectors $`\stackrel{}{a}^{(k)}`$ and $`\stackrel{}{b}^{(k,j)}`$ in the homogeneous relations (A.21) and (A.22) satisfied by $`𝐅_{𝒟,c}(x,t,\lambda )`$ being independent of $`x`$ and $`t`$, these relations are satisfied by $`(𝒳𝐅_{𝒟,c})(x,t,\lambda )`$ and $`(𝒯𝐅_{𝒟,c})(x,t,\lambda )`$ as well. For example, with the operator $`𝒳`$,
$$(𝒳(\lambda ,\stackrel{}{q})𝐅_{𝒟,c})(x,t,\lambda _k)\stackrel{}{a}^{(k)}=𝒳(\lambda _k,\stackrel{}{q})(𝐅_{𝒟,c}(x,t,\lambda _k)\stackrel{}{a}^{(k)})=𝒳(\lambda _k,\stackrel{}{q})\mathrm{𝟎}=\mathrm{𝟎},$$
(A.32)
for $`k=1,\mathrm{},M`$, and
$$(𝒳(\lambda ,\stackrel{}{q})𝐅_{𝒟,c})(x,t,\lambda _k^{})\stackrel{}{b}^{(k,j)}=𝒳(\lambda _k^{},\stackrel{}{q})(𝐅_{𝒟,c}(x,t,\lambda _k^{})\stackrel{}{b}^{(k,j)})=𝒳(\lambda _k^{},\stackrel{}{q})\mathrm{𝟎}=\mathrm{𝟎},$$
(A.33)
for $`k=1,\mathrm{},M`$ and $`j=1,\mathrm{},N`$. The argument is unchanged if $`𝒳`$ is replaced with $`𝒯`$.
Next, we examine the form of the matrix $`(𝒳(\lambda ,\stackrel{}{q})𝐅_{𝒟,c})(x,t,\lambda )`$. It is straightforward to see that
$$\begin{array}{ccc}\hfill (𝒳(\lambda ,\stackrel{}{q})𝐅_{𝒟,c})(x,t,\lambda )& =& \{\lambda ^M(2i[𝐄,𝐅^{(M1)}]c𝐔)+\hfill \\ & & \\ & & \underset{p=1}{\overset{M1}{}}\lambda ^p\left(_x𝐅^{(p)}+2i[𝐄,𝐅^{(p1)}]\mathrm{𝐔𝐅}^{(p)}\right)+\hfill \\ & & \\ & & (_x𝐅^{(0)}\mathrm{𝐔𝐅}^{(0)})\}\mathrm{exp}(2i(\lambda x+\lambda ^2t)𝐄).\hfill \end{array}$$
(A.34)
Now, as a consequence of the definition of the functions $`q_k(x,t)`$ in terms of the discrete data $`𝒟`$ and the symmetry property guaranteed by Proposition A.23, the leading term vanishes identically, that is,
$$2i[𝐄,𝐅^{(M1)}(x,t)]=c𝐔(\stackrel{}{q}).$$
(A.35)
This, along with the fact that $`𝒳(\lambda ,\stackrel{}{q})𝐅_{𝒟,c}(x,t,\lambda )`$ satisfies the homogeneous conditions (A.21) and (A.22) means that
$$𝒳(\lambda ,\stackrel{}{q})𝐅_{𝒟,c}(x,t,\lambda )\mathrm{\Lambda }_{N,M}^𝒟.$$
(A.36)
Not only that, but for matrices in $`\mathrm{\Lambda }_{N,M}^𝒟`$ the only way that the coefficient of $`\lambda ^M`$ can vanish is for the leading constant to vanish. Therefore, by the isomorphism between $`\mathrm{\Lambda }_{N,M}^𝒟`$ and $``$ via the leading constant, it follows that
$$𝒳(\lambda ,\stackrel{}{q})𝐅_{𝒟,c}(x,t,\lambda )=\mathrm{𝟎}.$$
(A.37)
We now consider the form of $`𝒯(\lambda ,\stackrel{}{q})𝐅_{𝒟,c}(x,t,\lambda )`$:
$$\begin{array}{ccc}\hfill 𝒯(\lambda ,\stackrel{}{q})𝐅_{𝒟,c}(x,t,\lambda )& =& \{\lambda ^{M+1}(2i[𝐄,𝐅^{(M1)}]c𝐔)+\hfill \\ & & \\ & & \lambda ^M\left(2i[𝐄,𝐅^{(M2)}]\mathrm{𝐔𝐅}^{(M1)}\frac{ic}{2}𝐕\right)+\hfill \\ & & \\ & & \underset{p=2}{\overset{M1}{}}\lambda ^p\left(_t𝐅^{(p)}+2i[𝐄,𝐅^{(p2)}]\mathrm{𝐔𝐅}^{(p1)}\frac{i}{2}\mathrm{𝐕𝐅}^{(p)}\right)+\hfill \\ & & \\ & & \lambda \left(_t𝐅^{(1)}\mathrm{𝐔𝐅}^{(0)}\frac{i}{2}\mathrm{𝐕𝐅}^{(1)}\right)+\hfill \\ & & \\ & & (_t𝐅^{(0)}\frac{i}{2}\mathrm{𝐕𝐅}^{(0)})\}\mathrm{exp}(2i(\lambda x+\lambda ^2t)𝐄).\hfill \end{array}$$
(A.38)
Once again, the definition of the functions $`q_k(x,t)`$ and the symmetry property of Proposition A.23 guarantee that the coefficient of $`\lambda ^{M+1}`$ vanishes. We shall now show that the coefficient of $`\lambda ^M`$ vanishes as well. Begin by writing $`𝐅^{(M1)}(x,t)`$ in the block form:
$$𝐅^{(M1)}(x,t)=\left[\begin{array}{cc}a^{(M1)}& \stackrel{}{b}^{(M1)T}\\ \stackrel{}{c}^{(M1)}& 𝐃^{(M1)}\end{array}\right].$$
(A.39)
We already know by definition of the functions $`q_k(x,t)`$ and Proposition A.23, that $`\stackrel{}{b}^{(M1)T}=\stackrel{}{q}^T/2i`$ and $`\stackrel{}{c}^{(M1)}=\stackrel{}{b}^{(M1)}`$. Making use of the fact that all of the terms in (A.34) vanish identically, we also have
$$_x𝐅^{(M1)}(x,t)+2i[𝐄,𝐅^{(M2)}(x,t)]𝐔(\stackrel{}{q})𝐅^{(M1)}(x,t)=\mathrm{𝟎}.$$
(A.40)
This implies that for the coefficient of $`\lambda ^M`$ in (A.38) to vanish, it will be enough to show that
$$c𝐕(\stackrel{}{q})=2i_x𝐅^{(M1)}(x,t).$$
(A.41)
From (A.4), it is clear (A.41) is satisfied in the off-diagonal blocks. To show that the diagonal blocks also vanish, we write out the diagonal blocks of (A.40):
$$\left[\begin{array}{cc}_xa^{(M1)}& \stackrel{}{0}^T\\ \stackrel{}{0}& _x𝐃^{(M1)}\end{array}\right]\left[\begin{array}{cc}\stackrel{}{q}^T\stackrel{}{c}^{(M1)}& \stackrel{}{0}^T\\ \stackrel{}{0}& \stackrel{}{q}^{}\stackrel{}{b}^{(M1)T}\end{array}\right]=\mathrm{𝟎}.$$
(A.42)
Eliminating the derivatives of $`a^{(M1)}`$ and $`𝐃^{(M1)}`$ between this equation and the diagonal blocks of (A.41) and comparing with the definition (A.4) of $`𝐕(\stackrel{}{q})`$, we finally see that (A.41) is satisfied identically. By similar arguments as we used above, it follows that $`𝒯(\lambda ,\stackrel{}{q})𝐅_{𝒟,c}(x,t,\lambda )`$ is also the zero element of $`\mathrm{\Lambda }_{N,M}^𝒟`$. This ends the proof of the proposition. $`\mathrm{}`$
We now return to the problem of interest, namely the algebraic construction of separable time-dependent potentials for the linear Schrödinger equation and of a large number of exact solutions to this linear equation. From the construction of the subspace $`\mathrm{\Lambda }_{N,M}^𝒟`$ we can extract a simpler construction of the quantities of immediate interest, and cast the whole procedure in the form of an algorithm. The key observation is that it is sufficient to build from given discrete data $`𝒟`$ only the first row of a matrix $`𝐅_{𝒟,c}(x,t,\lambda )`$ in the space $`\mathrm{\Lambda }_{N,M}^𝒟`$. This gives us both the functions $`q_k(x,t)`$ via the first row of the coefficient matrix $`𝐅^{(M1)}(x,t)`$ from which we find the potential $`V_0(x,t)`$ and also the image of the projection operator $`𝒫`$ that consists of solutions of the linear Schrödinger equation with this potential.
So we consider the first row of $`𝐅_{𝒟,c}(x,t,\lambda )`$ and impose the homogeneous linear constraints (A.21) and (A.22). Introducing
$$a(x,t,\lambda )=F_{11}(x,t,\lambda )=\left(\lambda ^M+\underset{p=0}{\overset{M1}{}}\lambda ^pa^{(p)}(x,t)\right)e^{2i(\lambda x+\lambda ^2t)},$$
(A.43)
and
$$\stackrel{}{b}(x,t,\lambda )=(F_{12}(x,t,\lambda ),\mathrm{},F_{1,N+1}(x,t,\lambda ))^T=\underset{p=0}{\overset{M1}{}}\lambda ^p\stackrel{}{b}^{(p)}(x,t),$$
(A.44)
the relations (A.21) and (A.22) take the simple form:
$$\begin{array}{ccc}\hfill a(x,t,\lambda _k)& =& \stackrel{}{g}^{(k)}\stackrel{}{b}(x,t,\lambda _k),\hfill \\ \hfill \stackrel{}{b}(x,t,\lambda _k^{})& =& a(x,t,\lambda _k^{})\stackrel{}{g}^{(k)},\hfill \end{array}$$
(A.45)
where $`k=1,\mathrm{},M`$. Note that without loss of generality, we are taking $`c=1`$. Written out in its entirety, this is a square linear system for the $`M(N+1)`$ unknowns, $`a^{(p)}(x,t)`$, and the $`N`$ elements of $`\stackrel{}{b}^{(p)}(x,t)`$ for $`p=0,\mathrm{},M1`$. The matrix of this system, and the right-hand side, are explicit functions of $`x`$ and $`t`$ through the exponential functions contributed by $`a(x,t,\lambda _k)`$ and $`a(x,t,\lambda _k^{})`$.
From the solution of this linear system, one computes the potential function as:
$$V_0(x,t)=4\underset{n=1}{\overset{N}{}}\left|b_n^{(M1)}(x,t)\right|^2.$$
(A.46)
Then, we see that $`a(x,t,\lambda )`$ and all the elements of $`\stackrel{}{b}(x,t,\lambda )`$ are solutions of the linear equation
$$i_tf+\frac{1}{2}_x^2fV_0(x,t)f=0,$$
(A.47)
for fixed but arbitrary $`\lambda `$. Being polynomial in $`\lambda `$, each element of $`\stackrel{}{b}(x,t,\lambda )`$ sweeps out an $`M`$-dimensional space of solutions as $`\lambda `$ varies. The solutions contained in $`a(x,t,\lambda )`$ are more interesting because the presence of the exponential means that all of these solutions for real $`\lambda `$ are linearly independent. This immediately gives an infinite-dimensional space of solutions to the linear Schrödinger equation.
In fact, the function $`a(x,t,\lambda )`$ contains an $`L^2()`$ basis of solutions of the Schrödinger equation as the parameter $`\lambda `$ is varied . In particular, the set
$$\{a(x,t,\lambda _1^{}),\mathrm{},a(x,t,\lambda _M^{}),a(x,t,\lambda ),\lambda \mathrm{real}\}$$
(A.48)
considered as functions of $`x`$ for fixed $`t`$, is complete. For real $`\lambda `$, set
$$\mathrm{\Psi }_\mathrm{d}(x,t,\lambda )\left(\pi \underset{k=1}{\overset{M}{}}|\lambda \lambda _k|^2\right)^{1/2}a(x,t,\lambda ).$$
(A.49)
The subscript “d” indicates solutions that superpose to form dispersive waves. For $`\lambda `$ and $`\eta `$ real we then have the inner products
$$\mathrm{\Psi }_\mathrm{d}(,t,\lambda ),\mathrm{\Psi }_\mathrm{d}(,t,\eta )=\delta (\lambda \eta ),$$
(A.50)
and for $`k=1,\mathrm{},M`$,
$$a(,t,\lambda _k^{}),\mathrm{\Psi }_\mathrm{d}(,t,\mu )=0.$$
(A.51)
Also, $`\mathrm{dim}\mathrm{span}\{a(x,t,\lambda _k^{}),k=1,\mathrm{},M\}=M`$ as functions of $`x`$ for fixed $`t`$. So, let $`\{\mathrm{\Psi }_{\mathrm{b},k}(x,t)\}`$ be any basis of $`\mathrm{span}\{a(x,t,\lambda _k^{}),k=1,\mathrm{},M\}`$ that is orthonormal with respect to the inner product (say obtained by the Gram-Schmidt procedure), so that
$$\mathrm{\Psi }_{\mathrm{b},j}(,t),\mathrm{\Psi }_{\mathrm{b},k}(,t)=\delta _{jk}.$$
(A.52)
The subscript “b” indicates solutions that are bound and have finite energy. Note that this basis remains orthonormal because the time evolution of these functions under (A.47) is unitary.
The completeness relation is generalized to $`L^2()`$ from that proved in as:
###### Proposition A.5
Let discrete data $`𝒟`$ be given and let $`t`$ be fixed. For all $`f(x)L^2()`$, we have the expansion
$$f(x)=_{\mathrm{}}^{\mathrm{}}f_\mathrm{d}(\lambda ,t)\mathrm{\Psi }_\mathrm{d}(x,t,\lambda )𝑑\lambda +\underset{k=1}{\overset{M}{}}f_{\mathrm{b},k}(t)\mathrm{\Psi }_{\mathrm{b},k}(x,t),$$
(A.53)
where the expansion coefficients are given by:
$$f_\mathrm{d}(\lambda ,t)=\mathrm{\Psi }_\mathrm{d}(,t,\lambda ),f(),f_{\mathrm{b},k}(t)=\mathrm{\Psi }_{\mathrm{b},k}(,t),f().$$
(A.54)
The orthogonality relations for the functions $`\mathrm{\Psi }_{\mathrm{b},k}(x,t)`$ and $`\mathrm{\Psi }_\mathrm{d}(x,t,\lambda )`$ are implied by this result. Note that if $`f=f(x,t)`$ satisfies (A.47) then the expansion coefficients are independent of $`t`$ and can be constructed from the initial data $`f(x,0)`$. Thus one solves the initial value problem for (A.47) in $`L^2()`$.
## Appendix Appendix B Dispersive Local Decay Estimates
Here, we establish several important properties of the unitary group $`e^{it}`$. We will consider even perturbations of the even two-soliton periodic potentials, so we will work in either the even or odd subspace of $`L^2()`$. For a given function $`f()`$ in $`L^2()`$, the operator $`𝒫_\mathrm{c}^{(\mathrm{e},\mathrm{o})}`$ is defined as the spectral projection onto the continuous part of the spectrum of $``$:
$$(𝒫_\mathrm{c}^{(\mathrm{e},\mathrm{o})}f)(x)_0^{\mathrm{}}\mathrm{\Psi }_\mathrm{d}^{(\mathrm{e},\mathrm{o})}(,0,\lambda ),f()\mathrm{\Psi }_\mathrm{d}^{(\mathrm{e},\mathrm{o})}(x,0,\lambda )𝑑\lambda .$$
(B.1)
As we will now see, the main difference between the even and odd cases is in the rate of dispersive decay, and the difference can be directly traced to the behavior of the dispersive eigenfunction $`\mathrm{\Psi }_\mathrm{d}^{(\mathrm{e},\mathrm{o})}(x,0,\lambda )`$ in the vicinity of $`\lambda =0`$. It is easy to see from the explicit formulas that the eigenfunctions are continuous in $`\lambda `$ at $`\lambda =0`$, and that the odd mode vanishes there:
$$\mathrm{\Psi }_\mathrm{d}^{(\mathrm{o})}(x,0,\lambda =0)=0,$$
(B.2)
while the even mode does not vanish, but is simply finite at $`\lambda =0`$. We say that the existence of a nontrivial eigenfunction at $`\lambda =0`$, as in the even case, indicates a zero-energy resonance of the system. The ubiquitous effect of a zero-energy resonance is to alter the rate of dispersive decay in the system. However, more dramatic effects can appear if under the influence of a perturbation, the zero-energy resonance is directly excited. This latter situation we refer to as a parametric zero-energy resonance. A system with a zero-energy resonance is “primed” to feel the effects of a parametric zero-energy resonance in the presence of an appropriate perturbation.
### B.1 Nonsingular local decay.
First, we will prove the nonsingular local decay estimate for the unitary group $`e^{it}`$.
###### Proposition B.1
Fix $`\sigma >5/2`$. There exist constants $`L^{(\mathrm{e},\mathrm{o})}>0`$ such that
$$^\sigma \left(e^{it}𝒫_\mathrm{c}^{(\mathrm{e})}f\right)()_2L^{(\mathrm{e})}t^{1/2}^\sigma f()_2,$$
(B.3)
and
$$^\sigma \left(e^{it}𝒫_\mathrm{c}^{(\mathrm{o})}f\right)()_2L^{(\mathrm{o})}t^{3/2}^\sigma f()_2,$$
(B.4)
for all $`fL^2()`$ for which the right hand side makes sense.
The proof is based on a sequence of intermediate results. First, from the simple chain of estimates:
$$^\sigma e^{it}𝒫_\mathrm{c}^{(\mathrm{e},\mathrm{o})}f()_2e^{it}𝒫_\mathrm{c}^{(\mathrm{e},\mathrm{o})}f()_2=𝒫_\mathrm{c}^{(\mathrm{e},\mathrm{o})}f()_2f()_2^\sigma f()_2,$$
(B.5)
we have
###### Lemma B.1
For all $`\sigma >0`$, we have the simple estimate
$$^\sigma e^{it}𝒫_\mathrm{c}^{(\mathrm{e},\mathrm{o})}f()_2^\sigma f()_2,$$
(B.6)
for all $`fL^2()`$ for which the right-hand side make sense.
We now want to refine the above uniform estimate to include a multiplicative factor of $`^\sigma f()_2`$ that decays in $`|t|`$. To this end, we fix $`t0`$ and observe that by the definition of the operator $``$,
$$x^\sigma \left(e^{it}𝒫_\mathrm{c}^{(\mathrm{e},\mathrm{o})}f\right)(x)=_{\mathrm{}}^{\mathrm{}}y^\sigma f(y)h(x,y;t)𝑑y,$$
(B.7)
where
$$h(x,y;t)x^\sigma y^\sigma _0^{\mathrm{}}\mathrm{\Psi }_\mathrm{d}^{(\mathrm{e},\mathrm{o})}(y,0,\lambda )^{}\mathrm{\Psi }_\mathrm{d}^{(\mathrm{e},\mathrm{o})}(x,0,\lambda )e^{2i\lambda ^2t}𝑑\lambda .$$
(B.8)
We note here that the integral in the definition of $`h(x,y;t)`$ is improper; the integrand is not absolutely integrable, and the integral from zero to infinity should be interpreted as the limit of the integral from zero to $`R`$ as $`R\mathrm{}`$. This limit exists as long as $`t0`$, and consequently the function $`h(x,y;t)`$ is well-defined for $`t0`$. The trouble with the function $`h(x,y;t)`$ at $`t=0`$ is not our concern here because we already have a uniform estimate that holds for all $`t`$, and in particular for $`t`$ near zero. Thus we will be thinking of $`t`$ as being large in what follows.
In any case, by Cauchy-Schwarz, we have
$$\left|x^\sigma \left(e^{it}𝒫_\mathrm{c}^{(\mathrm{e},\mathrm{o})}f\right)(x)\right|h(x,;t)_2^\sigma f()_2.$$
(B.9)
It follows that
$$^\sigma \left(e^{it}𝒫_\mathrm{c}^{(\mathrm{e},\mathrm{o})}f\right)()_2h(,;t)_2^\sigma f()_2,$$
(B.10)
an estimate that involves the Hilbert-Schmidt norm of the kernel $`h(x,y;t)`$ for each fixed $`t`$. The rest of our work will be to show $`h(x,y;t)`$ is in $`L^2(^2)`$ for each fixed $`t`$, with norm decaying in $`|t|`$.
First note that from the explicit formulas:
$$\mathrm{\Psi }_\mathrm{d}^{(\mathrm{o})}(x,0,\lambda )=\frac{2\lambda a^{(1)}(x,0)\mathrm{cos}(2\lambda x)2i(\lambda ^2+a^{(0)}(x,0))\mathrm{sin}(2\lambda x)}{\sqrt{2\pi (\lambda ^2+\rho _1^2)(\lambda ^2+\rho _2^2)}},$$
(B.11)
$$\mathrm{\Psi }_\mathrm{d}^{(\mathrm{e})}(x,0,\lambda )=\frac{2(\lambda ^2+a^{(0)}(x,0))\mathrm{cos}(2\lambda x)2i\lambda a^{(1)}(x,0)\mathrm{sin}(2\lambda x)}{\sqrt{2\pi (\lambda ^2+\rho _1^2)(\lambda ^2+\rho _2^2)}},$$
(B.12)
where $`a^{(0)}(x,t)`$ and $`a^{(1)}(x,t)`$ are bounded analytic functions of $`x`$, we obtain
###### Lemma B.2
Let the parameters $`\rho _1`$ and $`\rho _2`$ be fixed. The function defined by
$$q(\lambda )\mathrm{\Psi }_\mathrm{d}^{(\mathrm{e},\mathrm{o})}(x,0,\lambda )\mathrm{\Psi }_\mathrm{d}^{(\mathrm{e},\mathrm{o})}(y,0,\lambda )^{}$$
(B.13)
is in $`C^k(_+)`$ for all $`k0`$. In particular all derivatives with respect to $`\lambda `$ are uniformly bounded functions of $`\lambda `$. The norms $`q^{(k)}()_{\mathrm{}}`$ are homogeneous polynomials in $`|x|`$ and $`|y|`$ of degree $`k`$, with nonnegative coefficients that depend only on $`\rho _1`$ and $`\rho _2`$. Also, in the odd case, we have $`q(\lambda )=O(\lambda ^2)`$ for $`\lambda `$ near zero, while in the even case $`q(\lambda )=O(1)`$.
In showing that $`h(x,y;t)`$ is $`L^2(^2)`$ with norm decaying in $`t`$, we will find that the main contribution for large $`t`$ comes from the part of the integral near $`\lambda =0`$. To see this, we first separate the contributions near and away from zero. Let $`g_\mathrm{\Delta }(\lambda )`$ be a nonnegative “bump function”, infinitely differentiable for real $`\lambda >0`$, identically equal to $`1`$ for $`0\lambda \mathrm{\Delta }/2`$ and identically equal to zero for $`\lambda 3\mathrm{\Delta }/2`$. Let $`\stackrel{~}{g}_\mathrm{\Delta }(\lambda )1g_\mathrm{\Delta }(\lambda )`$. Then
$$h(x,y;t)=h_\mathrm{\Delta }(x,y;t)+\stackrel{~}{h}_\mathrm{\Delta }(x,y;t),$$
(B.14)
where
$$h_\mathrm{\Delta }(x,y;t)x^\sigma y^\sigma _0^{3\mathrm{\Delta }/2}\mathrm{\Psi }_\mathrm{d}^{(\mathrm{e},\mathrm{o})}(y,0,\lambda )^{}\mathrm{\Psi }_\mathrm{d}^{(\mathrm{e},\mathrm{o})}(x,0,\lambda )g_\mathrm{\Delta }(\lambda )e^{2i\lambda ^2t}𝑑\lambda ,$$
(B.15)
and
$$\stackrel{~}{h}_\mathrm{\Delta }(x,y;t)x^\sigma y^\sigma _{\mathrm{\Delta }/2}^{\mathrm{}}\mathrm{\Psi }_\mathrm{d}^{(\mathrm{e},\mathrm{o})}(y,0,\lambda )^{}\mathrm{\Psi }_\mathrm{d}^{(\mathrm{e},\mathrm{o})}(x,0,\lambda )\stackrel{~}{g}_\mathrm{\Delta }(\lambda )e^{2i\lambda ^2t}𝑑\lambda .$$
(B.16)
First, we will show that away from $`\lambda =0`$, we can get arbitrary decay in time.
###### Lemma B.3
Fix $`L>0`$. For some $`k2`$, suppose that $`f(\lambda )`$ is in $`C^n([L,\mathrm{}])`$ for all $`n=0,1,\mathrm{},k`$. Suppose that $`f(L)=f^{}(L)=\mathrm{}=f^{(k1)}(L)=0`$ and that the limit
$$\underset{R\mathrm{}}{lim}_L^Rf(\lambda )e^{2i\lambda ^2t}𝑑\lambda $$
(B.17)
exists for $`t0`$. Then
$$\left|\underset{R\mathrm{}}{lim}_L^Rf(\lambda )e^{2i\lambda ^2t}𝑑\lambda \right|\frac{1}{L4^k|t|^k}\underset{\lambda >L}{sup}\left|\lambda ^2(𝒜^kf)(\lambda )\right|,$$
(B.18)
where the operator $`𝒜`$ is defined by
$$(𝒜f)(\lambda )\frac{}{\lambda }\left(\frac{f(\lambda )}{\lambda }\right).$$
(B.19)
Proof: Integrating by parts $`k`$ times,
$$\begin{array}{ccc}\hfill \underset{R\mathrm{}}{lim}_L^Rf(\lambda )e^{2i\lambda ^2t}𝑑\lambda & =& \underset{R\mathrm{}}{lim}[\underset{n=0}{\overset{k1}{}}\left(\frac{i}{4t}\right)^{n+1}(1)^n\lambda ^1e^{2i\lambda ^2t}(𝒜^nf)(\lambda )|_{\lambda =L}^{\lambda =R}\hfill \\ & & \\ & & +\left(\frac{i}{4t}\right)^k_L^R(𝒜^kf)(\lambda )e^{2i\lambda ^2t}d\lambda ].\hfill \end{array}$$
(B.20)
The boundary terms at $`\lambda =L`$ vanish identically, and those at $`\lambda =R`$ tend to zero as $`R\mathrm{}`$. These facts prove the existence of the limit of the integral in the second line, and we find
$$\begin{array}{ccc}\hfill \left|\underset{R\mathrm{}}{lim}_L^Rf(\lambda )e^{2i\lambda ^2t}𝑑\lambda \right|& =& \frac{1}{4^k|t|^k}\left|\underset{R\mathrm{}}{lim}_L^R(𝒜^kf)(\lambda )e^{2i\lambda ^2t}𝑑\lambda \right|\hfill \\ & & \\ & & \frac{1}{4^k|t|^k}\underset{R\mathrm{}}{lim}_L^R\left|\lambda ^2(𝒜^kf)(\lambda )\right|\frac{d\lambda }{\lambda ^2}\hfill \\ & & \\ & & \frac{1}{L4^k|t|^k}\underset{\lambda >L}{sup}\left|\lambda ^2(𝒜^kf)(\lambda )\right|.\hfill \end{array}$$
(B.21)
The bound is finite for $`k2`$. $`\mathrm{}`$
We can now apply this result to estimate $`\stackrel{~}{h}_\mathrm{\Delta }(x,y;t)`$.
###### Lemma B.4
Fix an integer $`k2`$, and let $`\sigma >k+1/2`$. Then, the function $`\stackrel{~}{h}_\mathrm{\Delta }(x,y;t)`$ is in $`L^2(^2)`$ as a function of $`x`$ and $`y`$, with norm decaying as $`|t|^k`$.
Proof: We apply the above lemma with $`L=\mathrm{\Delta }/2`$ and $`f(\lambda )=\mathrm{\Psi }_\mathrm{d}^{(\mathrm{e},\mathrm{o})}(x,0,\lambda )\mathrm{\Psi }_\mathrm{d}^{(\mathrm{e},\mathrm{o})}(y,0,\lambda )^{}\stackrel{~}{g}_\mathrm{\Delta }(\lambda )`$. This gives the pointwise estimate
$$|\stackrel{~}{h}_\mathrm{\Delta }(x,y;t)|\frac{2x^\sigma y^\sigma }{\mathrm{\Delta }4^k|t|^k}\underset{\lambda >\mathrm{\Delta }/2}{sup}\left|\lambda ^2(𝒜^k\mathrm{\Psi }_\mathrm{d}^{(\mathrm{e},\mathrm{o})}(x,0,)\mathrm{\Psi }_\mathrm{d}^{(\mathrm{e},\mathrm{o})}(y,0,)^{}\stackrel{~}{g}_\mathrm{\Delta }())(\lambda )\right|.$$
(B.22)
The operator $`𝒜^k`$ acting on the right-hand side makes the supremum bound a polynomial in $`|x|`$ and $`|y|`$ of degree $`k`$. Therefore for $`\stackrel{~}{h}_\mathrm{\Delta }(x,y;t)`$ to lie in $`L^2(^2)`$ as a function of $`x`$ and $`y`$, it is sufficient to take $`\sigma >k+1/2`$. The claimed time decay of the $`L^2`$ norm is then obvious. Note that each derivative of $`\stackrel{~}{g}_\mathrm{\Delta }(\lambda )`$ contributes a factor of order $`O(\mathrm{\Delta }^1)`$, so the overall bound on the $`L^2`$ norm of $`\stackrel{~}{h}_\mathrm{\Delta }(x,y;t)`$ scales like $`\mathrm{\Delta }^{(k+1)}`$. $`\mathrm{}`$
Now, we move on to consider the part of $`h(x,y;t)`$ contributed by the neighborhood of $`\lambda =0`$. We again need some technical lemmas.
###### Lemma B.5
For all $`\mu `$,
$$\left|_0^\mu e^{2i\zeta ^2}𝑑\zeta \right|\sqrt{3}.$$
(B.23)
Proof: First, note that
$$\left|_0^\mu e^{2i\zeta ^2}𝑑\zeta \right|_0^\mu |d\zeta |=|\mu |.$$
(B.24)
This estimate is useful for bounded $`\mu `$. Suppose $`\mu >M>0`$. Then,
$$\left|_0^\mu e^{2i\zeta ^2}𝑑\zeta \right|M+\left|_M^\mu e^{2i\zeta ^2}𝑑\zeta \right|.$$
(B.25)
Changing variables to $`\tau =\zeta ^2`$ and integrating by parts, one finds
$$\left|_M^\mu e^{2i\zeta ^2}𝑑\zeta \right|=\left|\frac{ie^{2i\mu ^2}}{4\mu }\frac{ie^{2iM^2}}{4M}+\frac{i}{8}_{M^2}^{\mu ^2}e^{2i\tau }\tau ^{3/2}𝑑\tau \right|\frac{3}{4M}.$$
(B.26)
Therefore, for $`\mu >M>0`$, we have the estimate
$$\left|_0^\mu e^{2i\zeta ^2}𝑑\zeta \right|M+\frac{3}{4M}.$$
(B.27)
The right hand side takes its smallest value, $`\sqrt{3}`$, for $`M_{\mathrm{min}}=\sqrt{3}/2`$. Since for $`0<\mu M_{\mathrm{min}}`$, we have
$$\left|_0^\mu e^{2i\zeta ^2}𝑑\zeta \right||\mu |M_{\mathrm{min}}2M_{\mathrm{min}}=\sqrt{3},$$
(B.28)
the lemma is established uniformly for all positive $`\mu `$. By symmetry, the same estimate holds for $`\mu <0`$. $`\mathrm{}`$
###### Lemma B.6
Fix $`L>0`$ and suppose $`f(\lambda )`$ is twice continuously differentiable, with $`f(0)=f^{}(0)=0`$, and $`f(L)=f^{}(L)=0`$. Then
$$\left|_0^Lf(\lambda )e^{2i\lambda ^2t}𝑑\lambda \right|\frac{L\sqrt{3}}{4|t|^{3/2}}\underset{0<\lambda <L}{sup}\left|\frac{^2}{\lambda ^2}\left(\frac{f(\lambda )}{\lambda }\right)\right|.$$
(B.29)
Proof: Integrating by parts using the boundary conditions (evaluations at the lower boundary of $`\lambda =0`$ are interpreted in the sense of the limit $`\lambda 0`$, that is, from above), we have
$$_0^Lf(\lambda )e^{2i\lambda ^2t}𝑑\lambda =\frac{i}{4t}_0^L\frac{f(\lambda )}{\lambda }\frac{}{\lambda }\left(e^{2i\lambda ^2t}\right)𝑑\lambda =\frac{i}{4t}_0^L\frac{}{\lambda }\left(\frac{f(\lambda )}{\lambda }\right)e^{2i\lambda ^2t}𝑑\lambda .$$
(B.30)
Write
$$e^{2i\lambda ^2t}=\frac{}{\lambda }_0^\lambda e^{2i\sigma ^2t}𝑑\sigma ,$$
(B.31)
and integrate by parts again making use of the boundary conditions (with the same caveat as above), to find
$$_0^Lf(\lambda )e^{2i\lambda ^2t}𝑑\lambda =\frac{i}{4t}_0^L\frac{^2}{\lambda ^2}\left(\frac{f(\lambda )}{\lambda }\right)_0^\lambda e^{2i\sigma ^2t}𝑑\sigma 𝑑\lambda .$$
(B.32)
With a change of variables to $`\zeta =|t|^{1/2}\sigma `$, this becomes
$$_0^Lf(\lambda )e^{2i\lambda ^2t}𝑑\lambda =\frac{i}{4t|t|^{1/2}}_0^L\frac{^2}{\lambda ^2}\left(\frac{f(\lambda )}{\lambda }\right)_0^{|t|^{1/2}\lambda }e^{2i\zeta ^2}𝑑\zeta 𝑑\lambda .$$
(B.33)
Estimating the $`\lambda `$ integral in the obvious way using the uniform bound of the $`\zeta `$ integral by $`\sqrt{3}`$ establishes the claimed estimate. $`\mathrm{}`$
Without the vanishing boundary conditions at $`\lambda =0`$, one gets less decay in time.
###### Lemma B.7
Let $`f(\lambda )`$ be absolutely continuous $`0\lambda L`$, so that $`f^{}(\lambda )L^1([0,L])`$. Then
$$\left|_0^Lf(\lambda )e^{2i\lambda ^2t}𝑑\lambda \right|\left(|f(0)|+2_0^L|f^{}(\lambda )|𝑑\lambda \right)\frac{\sqrt{3}}{|t|^{1/2}},$$
(B.34)
an order $`O(|t|^{1/2})`$ bound.
Proof: Separate off the slow decay by writing
$$_0^Lf(\lambda )e^{2i\lambda ^2t}𝑑\lambda =f(0)_0^Le^{2i\lambda ^2t}𝑑\lambda +_0^L(f(\lambda )f(0))e^{2i\lambda ^2t}𝑑\lambda =I_A+I_B.$$
(B.35)
The first integral is easily transformed:
$$I_A=f(0)_0^Le^{2i\lambda ^2t}𝑑\lambda =\frac{f(0)}{t^{1/2}}_0^{Lt^{1/2}}e^{2i\zeta ^2}𝑑\zeta ,$$
(B.36)
and therefore easily uniformly estimated
$$|I_A|\frac{\sqrt{3}|f(0)|}{|t|^{1/2}}.$$
(B.37)
In the second integral, one integrates by parts to find
$$I_B=_0^Lf^{}(\lambda )_\lambda ^Le^{2i\mu ^2t}𝑑\mu 𝑑\lambda .$$
(B.38)
Therefore,
$$\begin{array}{ccc}\hfill |I_B|& & \underset{0<\lambda <L}{sup}\left|_\lambda ^Le^{2i\mu ^2t}𝑑\mu \right|_0^L|f^{}(\lambda )|𝑑\lambda \hfill \\ & & \\ & & \left(\left|_0^Le^{2i\mu ^2t}𝑑\mu \right|+\underset{0<\lambda <L}{sup}\left|_0^\lambda e^{2i\mu ^2t}𝑑\mu \right|\right)_0^L|f^{}(\lambda )|𝑑\lambda \hfill \\ & & \\ & & \frac{2\sqrt{3}}{|t|^{1/2}}_0^L|f^{}(\lambda )|𝑑\lambda .\hfill \end{array}$$
(B.39)
Combining the estimates for $`I_A`$ and $`I_B`$ establishes the claimed result. $`\mathrm{}`$
We now want to use these results to estimate $`h_\mathrm{\Delta }(x,y;t)`$. To do this, we want to apply Lemma B.6 or Lemma B.7 with $`f(\lambda )=\mathrm{\Psi }_\mathrm{d}^{(\mathrm{e},\mathrm{o})}(x,0,\lambda )\mathrm{\Psi }_\mathrm{d}^{(\mathrm{e},\mathrm{o})}(y,0,\lambda )^{}g_\mathrm{\Delta }(\lambda )`$. Now, from Lemma B.2, it is clear that the hypotheses of Lemma B.6 concerning the behavior of $`f`$ at $`\lambda =0`$ will only be satisfied in the odd case. Here, we obtain the following.
###### Lemma B.8
Consider the odd case, and let $`\sigma >5/2`$. Then $`h_\mathrm{\Delta }(x,y;t)`$ is in $`L^2(^2)`$ as a function of $`x`$ and $`y`$ with norm decaying like $`|t|^{3/2}`$.
Proof: We have the pointwise estimate
$$|h_\mathrm{\Delta }(x,y;t)|\frac{3\sqrt{3}\mathrm{\Delta }x^\sigma y^\sigma }{8|t|^{3/2}}\underset{0<\lambda <3\mathrm{\Delta }/2}{sup}\left|\frac{^2}{\lambda ^2}\left(\frac{\mathrm{\Psi }_\mathrm{d}^{(\mathrm{o})}(x,0,\lambda )\mathrm{\Psi }_\mathrm{d}^{(\mathrm{o})}(y,0,\lambda )^{}g_\mathrm{\Delta }(\lambda )}{\lambda }\right)\right|.$$
(B.40)
From Lemma B.2 we have that the right-hand side is a quadratic polynomial in $`|x|`$ and $`|y|`$. Therefore for $`h_\mathrm{\Delta }(x,y;t)`$ to be in $`L^2(^2)`$ as a function of $`x`$ and $`y`$ it is sufficient to take $`\sigma >5/2`$. The time decay of the $`L^2`$ norm is then obvious. Note that each derivative of $`g_\mathrm{\Delta }(\lambda )`$ contributes a factor that is $`O(\mathrm{\Delta }^1)`$ so the bound on the $`L^2`$ norm scales like $`\mathrm{\Delta }^1`$. $`\mathrm{}`$
In the even case, we have a zero-energy resonance, and this means that the integrand near $`\lambda =0`$ is not small enough to allow decay as rapid as in the odd case. In this case, we can only apply Lemma B.7 to find the following.
###### Lemma B.9
Consider the even case, and let $`\sigma >3/2`$. Then $`h_\mathrm{\Delta }(x,y;t)`$ is in $`L^2(^2)`$ as a function of $`x`$ and $`y`$, with norm decaying like $`|t|^{1/2}`$.
Proof: Using $`f(\lambda )=\mathrm{\Psi }_\mathrm{d}^{(\mathrm{e})}(x,0,\lambda )\mathrm{\Psi }_\mathrm{d}^{(\mathrm{e})}(y,0,\lambda )^{}g_\mathrm{\Delta }(\lambda )`$ and $`L=3\mathrm{\Delta }/2`$ in Lemma B.7, we have the pointwise estimate
$$|h_\mathrm{\Delta }(x,y;t)|\frac{\sqrt{3}x^\sigma y^\sigma }{|t|^{1/2}}\left(|f(0)|+2_0^{3\mathrm{\Delta }/2}|f^{}(\lambda )|𝑑\lambda \right).$$
(B.41)
Since the derivative with respect to $`\lambda `$ results in at most linear growth in $`x`$ and $`y`$, taking $`\sigma >3/2`$ is sufficient to ensure that $`h_\mathrm{\Delta }(x,y;t)`$ is in $`L^2(^2)`$ as a function of $`x`$ and $`y`$. Clearly, for large $`t`$, the $`L^2`$ norm is $`O(|t|^{1/2})`$. Note that the estimate is also $`O(\mathrm{\Delta }^1)`$ due to differentiation of the bump function $`g_\mathrm{\Delta }(\lambda )`$. $`\mathrm{}`$
In both odd and even cases, the contribution of $`h_\mathrm{\Delta }(x,y;t)`$ to the $`L^2`$ norm of $`h(x,y;t)`$ dominates for large time that of $`\stackrel{~}{h}_\mathrm{\Delta }(x,y;t)`$, for which we had arbitrary decay. According to Lemma B.4, for $`\sigma >5/2`$ this latter decay is at least as fast as $`|t|^2`$. These results imply the following.
###### Lemma B.10
Fix $`\sigma >5/2`$. Then, for $`t`$ sufficiently large, we have the estimates:
$$^\sigma \left(e^{it}𝒫_\mathrm{c}^{(\mathrm{e})}f\right)()_2\frac{K^{(\mathrm{e})}}{|t|^{1/2}}^\sigma f()_2,$$
(B.42)
and
$$^\sigma \left(e^{it}𝒫_\mathrm{c}^{(\mathrm{o})}f\right)()_2\frac{K^{(\mathrm{o})}}{|t|^{3/2}}^\sigma f()_2,$$
(B.43)
where $`K^{(\mathrm{e})}`$ and $`K^{(\mathrm{o})}`$ are some positive constants.
This result, taken together with the elementary time-independent bound established in Lemma B.1 completes the proof of Proposition B.1.
### B.2 Singular local decay.
Now we prove the singular local decay estimate for the unitary group $`e^{it}`$.
###### Proposition B.2
Let $`|\mu |\mu _{\mathrm{min}}>0`$. Fix $`\sigma >7/2`$. Let $`t=\kappa r`$ with $`r0`$ and $`\kappa =\pm 1`$. Then, there exist constants $`M^{(\mathrm{e},\mathrm{o})}>0`$, such that
$$^\sigma \underset{\delta 0}{lim}((2\mu 2i\kappa \delta )^1e^{it}𝒫_\mathrm{c}^{(\mathrm{e})}f)()_2M^{(\mathrm{e})}r^{1/2}^\sigma f()_2,$$
(B.44)
and
$$^\sigma \underset{\delta 0}{lim}((2\mu 2i\kappa \delta )^1e^{it}𝒫_\mathrm{c}^{(\mathrm{o})}f)()_2M^{(\mathrm{o})}r^{3/2}^\sigma f()_2.$$
(B.45)
The constants $`M^{(\mathrm{e},\mathrm{o})}`$ depend only on $`\mu _{\mathrm{min}}`$, so the bounds are uniform for large $`|\mu |`$.
The proof of this proposition begins with a representation similar to (B.7),
$$x^\sigma \underset{\delta 0}{lim}\left((2\mu 2i\kappa \delta )^1e^{it}𝒫_\mathrm{c}^{(\mathrm{e},\mathrm{o})}f\right)(x)=_{\mathrm{}}^{\mathrm{}}y^\sigma f(y)k(x,y;t)𝑑y,$$
(B.46)
where
$$k(x,y;t)=x^\sigma y^\sigma \underset{\delta 0}{lim}_0^{\mathrm{}}\frac{\mathrm{\Psi }_\mathrm{d}^{(\mathrm{e},\mathrm{o})}(x,0,\lambda )\mathrm{\Psi }_\mathrm{d}^{(\mathrm{e},\mathrm{o})}(y,0,\lambda )^{}}{2\lambda ^22\mu 2i\kappa \delta }e^{2i\lambda ^2t}𝑑\lambda .$$
(B.47)
Perhaps despite appearances, the kernel $`k(x,y;t)`$ is somewhat more amenable to analysis than the kernel $`h(x,y;t)`$ that appeared in the nonsingular case. This is because for each finite $`\delta `$ the integrand is absolutely integrable as a consequence of the uniform boundedness in $`\lambda `$ of $`\mathrm{\Psi }_\mathrm{d}^{(\mathrm{e},\mathrm{o})}(x,0,\lambda )`$ as guaranteed by Lemma B.2 and the large $`\lambda `$ behavior of the denominator.
Again, the goal is to show that the kernel $`k(x,y;t)`$ is in $`L^2(^2)`$ as a function of $`x`$ and $`y`$, with norm that is decaying in time, although in this case we will only obtain the decay for $`t`$ of a particular sign. First, we show that the $`L^2`$ norm exists and is finite near $`t=0`$.
###### Lemma B.11
Fix $`\sigma >3/2`$, $`\mu `$ with $`|\mu |\mu _{\mathrm{min}}>0`$, and $`t`$ with $`|t|<T`$. Then there exist constants $`C^{(\mathrm{e},\mathrm{o})}>0`$ depending on $`\mu _{\mathrm{min}}`$ and $`T`$ such that
$$k(,;t)_2C^{(\mathrm{e},\mathrm{o})}.$$
(B.48)
Since the bounds only depend on $`\mu `$ via $`\mu _{\mathrm{min}}`$, they are uniform for large $`|\mu |`$.
Proof: Begin by setting $`f(\lambda )=\mathrm{\Psi }_\mathrm{d}^{(\mathrm{e},\mathrm{o})}(x,0,\lambda )\mathrm{\Psi }_\mathrm{d}^{(\mathrm{e},\mathrm{o})}(y,0,\lambda )^{}e^{2i\lambda ^2t}`$. First we consider $`\mu \mu _{\mathrm{min}}<0`$, in which case we have
$$k(x,y;t)=\frac{x^\sigma y^\sigma }{2}_0^{\mathrm{}}\frac{f(\lambda )d\lambda }{\lambda ^2\mu },$$
(B.49)
since there is no singularity for $`\mu <0`$. We immediately get the pointwise estimate
$$\begin{array}{ccc}\hfill |k(x,y;t)|& & \frac{x^\sigma y^\sigma }{2}\underset{\lambda >0}{sup}|f(\lambda )|_0^{\mathrm{}}\frac{d\lambda }{\lambda ^2\mu }\hfill \\ & & \\ & =& \frac{\pi x^\sigma y^\sigma }{4\sqrt{\mu }}\underset{\lambda >0}{sup}|f(\lambda )|\hfill \\ & & \\ & & \frac{\pi x^\sigma y^\sigma }{4\sqrt{\mu _{\mathrm{min}}}}\underset{\lambda >0}{sup}|f(\lambda )|.\hfill \end{array}$$
(B.50)
Now we consider $`\mu \mu _{\mathrm{min}}>0`$. Pick some positive $`G`$ less than $`\sqrt{\mu _{\mathrm{min}}}`$. Then
$$k(x,y;t)=\frac{x^\sigma y^\sigma }{2}(I_{\mathrm{ns}}+I_\mathrm{s}),$$
(B.51)
where
$$\begin{array}{ccc}\hfill I_{\mathrm{ns}}& & _0^{\sqrt{\mu }G}\frac{f(\lambda )d\lambda }{\lambda ^2\mu }+_{\sqrt{\mu }+G}^{\mathrm{}}\frac{f(\lambda )d\lambda }{\lambda ^2\mu },\hfill \\ & & \\ \hfill I_\mathrm{s}& & \underset{\delta 0}{lim}_{\sqrt{\mu }G}^{\sqrt{\mu }+G}\frac{f(\lambda )d\lambda }{(\lambda +\sqrt{\mu +i\kappa \delta })(\lambda \sqrt{\mu +i\kappa \delta })},\hfill \end{array}$$
(B.52)
where the principal branch of the square root is understood, so that the square root is nearly a positive number for $`\delta `$ small. It is easy to find
$$\begin{array}{ccc}\hfill |I_{\mathrm{ns}}|& & \underset{\lambda >0}{sup}|f(\lambda )|\left(_0^{\sqrt{\mu }G}\frac{d\lambda }{\mu \lambda ^2}+_{\sqrt{\mu }+G}^{\mathrm{}}\frac{d\lambda }{\lambda ^2\mu }\right)\hfill \\ & & \\ & & \frac{sup_{\lambda >0}|f(\lambda )|}{\sqrt{\mu }}\left(\mathrm{arctanh}\left(\frac{\sqrt{\mu }}{\sqrt{\mu }+G}\right)+\mathrm{arctanh}\left(\frac{\sqrt{\mu }G}{\sqrt{\mu }}\right)\right)\hfill \\ & & \\ & & \underset{\lambda >0}{sup}|f(\lambda )|\underset{\mu >\mu _{\mathrm{min}}}{sup}\left(\frac{1}{\sqrt{\mu }}\left(\mathrm{arctanh}\left(\frac{\sqrt{\mu }}{\sqrt{\mu }+G}\right)+\mathrm{arctanh}\left(\frac{\sqrt{\mu }G}{\sqrt{\mu }}\right)\right)\right).\hfill \end{array}$$
(B.53)
For the singular part, we find
$$\begin{array}{ccc}\hfill I_\mathrm{s}& =& \underset{\delta 0}{lim}\frac{f(\sqrt{\mu })}{\sqrt{\mu }+\sqrt{\mu +i\kappa \delta }}_{\sqrt{\mu }G}^{\sqrt{\mu }+G}\frac{d\lambda }{\lambda \sqrt{\mu +i\kappa \delta }}+_{\sqrt{\mu }G}^{\sqrt{\mu }+G}\left(\frac{f(\lambda )}{\lambda +\sqrt{\mu }}\frac{f(\sqrt{\mu })}{2\sqrt{\mu }}\right)\frac{d\lambda }{\lambda \sqrt{\mu }}\hfill \\ & & \\ & =& \frac{i\pi \kappa f(\sqrt{\mu })}{2\sqrt{\mu }}+_{\sqrt{\mu }G}^{\sqrt{\mu }+G}\left(\frac{f(\lambda )}{\lambda +\sqrt{\mu }}\frac{f(\sqrt{\mu })}{2\sqrt{\mu }}\right)\frac{d\lambda }{\lambda \sqrt{\mu }}.\hfill \end{array}$$
(B.54)
Therefore,
$$\begin{array}{ccc}\hfill |I_\mathrm{s}|& & \frac{\pi |f(\sqrt{\mu })|}{2\sqrt{\mu }}+2G\underset{|\lambda \sqrt{\mu }|<G}{sup}\left|\left(\frac{f(\lambda )}{\lambda +\sqrt{\mu }}\frac{f(\sqrt{\mu })}{2\sqrt{\mu }}\right)\frac{1}{\lambda \sqrt{\mu }}\right|\hfill \\ & & \\ & & \frac{\pi |f(\sqrt{\mu })|}{2\sqrt{\mu }}+2G\underset{|\lambda \sqrt{\mu }|<G}{sup}\left|\frac{}{\lambda }\left(\frac{f(\lambda )}{\lambda +\sqrt{\mu }}\right)\right|.\hfill \end{array}$$
(B.55)
Again, the bounds are uniform in $`\mu `$ for large $`\mu `$.
Now, we simply note that the pointwise bounds for $`I_{\mathrm{ns}}`$ and $`I_\mathrm{s}`$ are themselves bounded by functions of $`x`$, $`y`$, and $`t`$ that grow linearly at worst, as a consequence of differentiation of $`f(\lambda )`$ with respect to $`\lambda `$ (c.f. Lemma B.2 for the growth in $`x`$ and $`y`$, while the growth in $`t`$ comes from the factor $`e^{2i\lambda ^2t}`$). Thus, to have $`k(x,y;t)L^2(^2)`$ as a function of $`x`$ and $`y`$, it is sufficient to take $`\sigma >3/2`$, and then the norm will be bounded by a linear function of $`|t|`$, and therefore uniformly for $`|t|<T`$. The bound is also uniform in $`\mu `$ for $`|\mu |\mu _{\mathrm{min}}>0`$. $`\mathrm{}`$
We now note that proving the decay for large $`r=|t|`$ for $`\mu \mu _{\mathrm{min}}<0`$ amounts to recalling the nonsingular local decay estimate. The integral is not really singular:
$$\underset{\delta 0}{lim}_0^{\mathrm{}}\frac{\mathrm{\Psi }_\mathrm{d}^{(\mathrm{e},\mathrm{o})}(x,0,\lambda )\mathrm{\Psi }_\mathrm{d}^{(\mathrm{e},\mathrm{o})}(y,0,\lambda )}{\lambda ^2\mu i\kappa \delta }e^{2i\lambda ^2t}𝑑\lambda =_0^{\mathrm{}}\frac{\mathrm{\Psi }_\mathrm{d}^{(\mathrm{e},\mathrm{o})}(x,0,\lambda )\mathrm{\Psi }_\mathrm{d}^{(\mathrm{e},\mathrm{o})}(y,0,\lambda )}{\lambda ^2\mu }e^{2i\lambda ^2t}𝑑\lambda .$$
(B.56)
Using the same arguments as used to prove the nonsingular local decay estimate one gets a pointwise bound for this integral that is at most quadratically growing in $`x`$ and $`y`$ and decaying like $`|t|^{1/2}`$ in the even case and $`|t|^{3/2}`$ in the odd case. Since the estimates involve up to two derivatives of the quotient in the integrand, the bounds will be uniform in $`\mu `$ for $`\mu \mu _{\mathrm{min}}<0`$.
To prove the decay for large $`r=|t|`$ in the truly singular case when $`\mu \mu _{\mathrm{min}}>0`$, we split $`k(x,y;t)`$ into three parts. Let $`g_\mathrm{\Delta }(\lambda )`$ and $`\stackrel{~}{g}_\mathrm{\Delta }(\lambda )`$ be as before, and introduce the new “bump” functions $`g_G(\lambda )`$ and $`\stackrel{~}{g}_G(\lambda )=1g_G(\lambda )`$, both infinitely differentiable and nonnegative, with $`g_G(\lambda )`$ identically equal to zero outside of the interval $`(\sqrt{\mu }3G/2,\sqrt{\mu }+3G/2)`$ and identically equal to one inside of the interval $`(\sqrt{\mu }G/2,\sqrt{\mu }+G/2)`$. Set
$$k(x,y;t)=\frac{x^\sigma y^\sigma }{2}\left(I_0+I_\mu +\stackrel{~}{I}\right),$$
(B.57)
where
$$I_0_0^{3\mathrm{\Delta }/2}\frac{\mathrm{\Psi }_\mathrm{d}^{(\mathrm{e},\mathrm{o})}(x,0,\lambda )\mathrm{\Psi }_\mathrm{d}^{(\mathrm{e},\mathrm{o})}(y,0,\lambda )^{}}{\lambda ^2\mu }g_\mathrm{\Delta }(\lambda )e^{2i\lambda ^2t}𝑑\lambda ,$$
(B.58)
$$I_\mu \underset{\delta 0}{lim}_{\sqrt{\mu }3G/2}^{\sqrt{\mu }+3G/2}\frac{\mathrm{\Psi }_\mathrm{d}^{(\mathrm{e},\mathrm{o})}(x,0,\lambda )\mathrm{\Psi }_\mathrm{d}^{(\mathrm{e},\mathrm{o})}(y,0,\lambda )^{}}{\lambda ^2\mu i\kappa \delta }g_G(\lambda )e^{2i\lambda ^2t}𝑑\lambda ,$$
(B.59)
and
$$\begin{array}{ccc}\hfill \stackrel{~}{I}& & _{\mathrm{\Delta }/2}^{\sqrt{\mu }G/2}\frac{\mathrm{\Psi }_\mathrm{d}^{(\mathrm{e},\mathrm{o})}(x,0,\lambda )\mathrm{\Psi }_\mathrm{d}^{(\mathrm{e},\mathrm{o})}(y,0,\lambda )^{}}{\lambda ^2\mu }\stackrel{~}{g}_\mathrm{\Delta }(\lambda )\stackrel{~}{g}_G(\lambda )e^{2i\lambda ^2t}𝑑\lambda \hfill \\ & & \\ & & +_{\sqrt{\mu }+G/2}^{\mathrm{}}\frac{\mathrm{\Psi }_\mathrm{d}^{(\mathrm{e},\mathrm{o})}(x,0,\lambda )\mathrm{\Psi }_\mathrm{d}^{(\mathrm{e},\mathrm{o})}(y,0,\lambda )^{}}{\lambda ^2\mu }\stackrel{~}{g}_G(\lambda )e^{2i\lambda ^2t}𝑑\lambda .\hfill \end{array}$$
(B.60)
Note that in keeping the contributions near zero and near $`\mu `$ distinct, we are assuming that $`3\mathrm{\Delta }/2<\sqrt{\mu _{\mathrm{min}}}G`$. The analysis of $`I_0`$ and $`\stackrel{~}{I}`$ proceeds exactly as in the proof of the nonsingular local decay estimate. The results are almost identical. For $`\stackrel{~}{I}`$ one can integrate by parts as many times as one likes, and therefore one gets a pointwise estimate with arbitrary decay in time of order $`O(|t|^k)`$ for $`k2`$, but at the cost of polynomial growth in $`x`$ and $`y`$ of degree $`k`$. For $`I_0`$, one gets a pointwise estimate that decays like $`O(|t|^{1/2})`$ in the even case and $`O(|t|^{3/2})`$ in the odd case, at the cost of quadratic growth in $`x`$ and $`y`$.
The estimates of $`I_0`$ and $`\stackrel{~}{I}`$ are uniform for large $`\mu `$. The pointwise bounds for $`\stackrel{~}{I}`$ involve supremum bounds over the range of integration of the quantity
$$\lambda ^2\left(𝒜^k\frac{f()}{()^2\mu }\right)(\lambda ),$$
(B.61)
with $`f(\lambda )`$ given by $`\mathrm{\Psi }_\mathrm{d}^{(\mathrm{e},\mathrm{o})}(x,0,\lambda )\mathrm{\Psi }_\mathrm{d}^{(\mathrm{e},\mathrm{o})}(y,0,\lambda )^{}\stackrel{~}{g}_G(\lambda )`$ for $`\lambda >\sqrt{\mu }+G/2`$ and with $`f(\lambda )`$ given by $`\mathrm{\Psi }_\mathrm{d}^{(\mathrm{e},\mathrm{o})}(x,0,\lambda )\mathrm{\Psi }_\mathrm{d}^{(\mathrm{e},\mathrm{o})}(y,0,\lambda )^{}\stackrel{~}{g}_\mathrm{\Delta }(\lambda )\stackrel{~}{g}_G(\lambda )`$ for $`\mathrm{\Delta }/2<\lambda <\sqrt{\mu }G/2`$. In particular, we will need these bounds for $`k=2`$, in which case
$$\lambda ^2\left(𝒜^k\frac{f()}{()^2\mu }\right)(\lambda )=\left(\frac{15\lambda ^410\mu \lambda ^2+3\mu }{\lambda ^2(\lambda ^2\mu )^3}\right)f(\lambda )+\left(\frac{7\lambda ^2+3\mu }{\lambda (\lambda ^2\mu )^2}\right)f^{}(\lambda )+\frac{f^{\prime \prime }(\lambda )}{\lambda ^2\mu }.$$
(B.62)
For the part of $`\stackrel{~}{I}`$ involving $`\lambda >\sqrt{\mu }+G/2`$, it is easy to check that the three coefficients above are monotonic functions of $`\lambda `$ for $`\lambda >\sqrt{\mu }`$ that decay for large $`\lambda `$ with $`\mu `$ fixed. Therefore each coefficient is bounded by its magnitude at the lower endpoint $`\lambda =\sqrt{\mu }+G/2`$. With $`G`$ held fixed, these bounds are then seen to be decaying functions of $`\mu `$. For the part of $`\stackrel{~}{I}`$ involving $`\lambda (\mathrm{\Delta }/2,\sqrt{\mu }G/2)`$, it is easy to see that the coefficients blow up at both endpoints. Therefore, for $`\mathrm{\Delta }`$ and $`G`$ sufficiently small but independent of $`\mu `$, the coefficients will be bounded by the maximum of their values at $`\lambda =\mathrm{\Delta }/2`$ and $`\lambda =\sqrt{\mu }G/2`$. Again, holding $`\mathrm{\Delta }`$ and $`G`$ fixed, one sees that the bounds are uniform for large $`\mu `$. This direct argument shows that, at least for $`k=2`$, the pointwise bound for $`\stackrel{~}{I}`$ is uniform in $`\mu `$. Establishing the uniformity of the pointwise estimate for $`I_0`$ is easier; the denominator $`\lambda ^2\mu `$ plays no essential role for $`\lambda <3\mathrm{\Delta }/2`$ for $`\mu `$ sufficiently large.
The new term that must be handled differently is $`I_\mu `$.
###### Lemma B.12
For all $`k2`$, the integral $`I_\mu `$ satisfies the pointwise estimate
$$|I_\mu |\frac{2}{k1}\frac{P_k(x,y)}{|t|^{k1}},$$
(B.63)
where $`P_k(x,y)`$ is a polynomial in $`|x|`$ and $`|y|`$ of degree $`k`$ with positive coefficients that are uniform in $`\mu `$.
Proof: Consider first $`t>0`$. Then, the quantity to estimate is
$$\begin{array}{ccc}\hfill I_\mu & =& \underset{\delta 0}{lim}_{\sqrt{\mu }3G/2}^{\sqrt{\mu }+3G/2}\frac{\mathrm{\Psi }_\mathrm{d}^{(\mathrm{e},\mathrm{o})}(x,0,\lambda )\mathrm{\Psi }_\mathrm{d}^{(\mathrm{e},\mathrm{o})}(y,0,\lambda )^{}}{\lambda ^2\mu i\delta }g_G(\lambda )e^{2i\lambda ^2t}𝑑\lambda \hfill \\ & & \\ & =& 2ie^{2i\mu t}\underset{\delta 0}{lim}e^{2\delta t}_{\sqrt{\mu }3G/2}^{\sqrt{\mu }+3G/2}\mathrm{\Psi }_\mathrm{d}^{(\mathrm{e},\mathrm{o})}(x,0,\lambda )\mathrm{\Psi }_\mathrm{d}^{(\mathrm{e},\mathrm{o})}(y,0,\lambda )^{}g_G(\lambda )_t^{\mathrm{}}e^{2i(\lambda ^2\mu i\delta )s}𝑑s𝑑\lambda \hfill \\ & & \\ & =& 2ie^{2i\mu t}_t^{\mathrm{}}e^{2i\mu s}\left[_{\sqrt{\mu }3G/2}^{\sqrt{\mu }+3G/2}\mathrm{\Psi }_\mathrm{d}^{(\mathrm{e},\mathrm{o})}(x,0,\lambda )\mathrm{\Psi }_\mathrm{d}^{(\mathrm{e},\mathrm{o})}(y,0,\lambda )^{}g_G(\lambda )e^{2i\lambda ^2s}𝑑\lambda \right]𝑑s.\hfill \end{array}$$
(B.64)
Now, with $`g_G(\lambda )`$ vanishing to all orders at the integration endpoints, it is possible to bound the integral in square brackets by iterated integration by parts. The bound is $`O(|s|^k)`$ and grows in $`x`$ and $`y`$ like a polynomial $`P_k(x,y)`$ of degree $`k`$. This bound is uniform in $`\mu `$, since the only place $`\mu `$ appears is in the range of integration over which bounds are required, and from Lemma B.2 we know that these bounds are uniform for all $`\lambda `$. Therefore we have
$$|I_\mu |2_t^{\mathrm{}}P_k(x,y)|s|^k𝑑s=\frac{2}{k1}\frac{P_k(x,y)}{|t|^{k1}},$$
(B.65)
which establishes the lemma for $`t>0`$. For $`t<0`$, one gets an integral from $`\mathrm{}`$ to $`t`$ in the second step above, and ultimately obtains the same bound. $`\mathrm{}`$
Finally, we put the pieces together to complete the proof of Proposition B.2. For the odd case, we want decay of order $`O(|t|^{3/2})`$. For $`x^\sigma y^\sigma I_0`$ to be in $`L^2(^2)`$ with this decay rate, we need $`\sigma >5/2`$. With this bound on $`\sigma `$, we can get $`x^\sigma y^\sigma \stackrel{~}{I}`$ being in $`L^2(^2)`$ with decay bounded by $`O(|t|^2)=o(|t|^{3/2})`$, but no better. Finally, for $`x^\sigma y^\sigma I_\mu `$ to be in $`L^2(^2)`$ with decay $`O(|t|^2)`$ we now see that we need to localize a bit more in space by taking $`\sigma >7/2`$. Combining these large time estimates with the finite time bound of Lemma B.11 establishes the proposition in the odd case. Similar arguments for the even case give an $`L^2(^2)`$ norm that decays like $`O(|t|^{1/2})`$ for $`\sigma >7/2`$. This finishes the proof of Proposition B.2.
Remark: Evidently, the singular decay estimates blow up when $`\mu _{\mathrm{min}}`$ approaches zero. This is an essential phenemenon in both the odd and even cases. This is best seen by considering the singular integral for the case $`\mu =0`$:
$$_0^{\mathrm{}}\frac{\mathrm{\Psi }_\mathrm{d}^{(\mathrm{e},\mathrm{o})}(x,0,\lambda )\mathrm{\Psi }_\mathrm{d}^{(\mathrm{e},\mathrm{o})}(y,0,\lambda )^{}}{\lambda ^2i\kappa \delta }e^{2i\lambda ^2t}𝑑\lambda .$$
(B.66)
This integral blows up for all $`x`$, $`y`$, and $`t`$, as $`\delta `$ tends to zero in the even case. In the odd case there is sufficient vanishing at $`\lambda =0`$ for the limit of $`\delta 0`$ to exist for all $`x`$, $`y`$, and $`t`$, but the limit only decays in $`t`$ like $`|t|^{1/2}`$. Thus it is not possible for estimates of the form derived for $`|\mu |\mu _{\mathrm{min}}>0`$ to hold uniformly in any neighborhood of $`\mu =0`$.$`\mathrm{}`$ |
warning/0002/hep-th0002062.html | ar5iv | text | # Constructing the light-front QCD Hamiltonian
## 1. Introduction
The search for ways to solve problems in the quantum field theory with a large coupling constant, has long been significant. Calculations on space-time lattices are now often used for this purpose in QCD. Essential results have been thus obtained. Nevertheless, these calculations are very laborious and have a low accuracy; in addition, it is generally difficult to estimate the calculation error theoretically. Therefore, it is interesting to seek other possible approaches to this problem. Even limited progress in this direction would allow comparing the results obtained by different methods.
Long before the advent of QCD, studying the pion-nucleon interaction had been attempted by solving the Schrodinger equation in the Lorentz coordinate system in the framework of the quantum theory of pion and nucleon fields. The states were described by the method previously found by V. A. Fock in terms of vectors in the space that now bears his name. In this case, the mathematical vacuum of the Fock space coincided with the free theory vacuum. This approach, now known as the Tamm-Dancoff method , did not lead to success. The primary reason was the complexity of the physical vacuum state, which did not coincide with the mathematical vacuum. Without a description of the physical vacuum, it was impossible to investigate any other states. If ultraviolet (UV) and infrared (IR) cutoffs are introduced to make the number of degrees of freedom finite, it would be possible, in principle, to represent the physical vacuum as a vector of the Fock space with the free theory vacuum. However, such a representation proves extremely complicated because it is necessary to provide the translation invariance of the physical vacuum and to satisfy the ”cluster decomposition property” for vacuum expectation values. For these reasons, the Schrodinger equation in the Lorentz coordinates, where the evolution occurs in the conventional time, is hardly applicable to any calculations in the quantum field theory with a large coupling constant.
As early as 1949, Dirac suggested a method for avoiding the difficulties related to the description of the physical vacuum state . He proposed using the light-front coordinates $`x^\pm =(x^0\pm x^3)/\sqrt{2}`$, $`x^1`$, $`x^2`$, where $`x^0`$, $`x^1`$, $`x^2`$, $`x^3`$ are the Lorentz coordinates, and treating $`x^+`$ as time. In this approach, a theory is canonically quantized on the surface $`x^+=const`$, and the generator $`P_+`$ of the shift along the $`x^+`$ axis plays the role of the Hamiltonian $`H`$. In addition, the generator of the shift along the $`x^{}`$ axis, i.e., the momentum operator $`P_{}`$, does not shift the surface $`x^+=const`$, where the quantization is performed, and is kinematic, according to the Dirac terminology (in contrast to the dynamic generator $`P_+`$). Therefore, the structure of the operator $`P_{}`$ in a theory with interaction is the same as in a theory without interaction, i.e., the operator $`P_{}`$ is always quadratic in the creation and annihilation operators $`a_{n}^{}{}_{}{}^{+}(p_{},p_{})`$ and $`a_n(p_{},p_{})`$ and, as a rule, has the form (after normal ordering)
$$\begin{array}{c}\text{ }P_{}=d^2p_{}\underset{0}{\overset{\mathrm{}}{}}𝑑p_{}p_{}\underset{n}{}a_{n}^{}{}_{}{}^{+}(p_{},p_{})a_n(p_{},p_{}),\text{ }\hfill \end{array}$$
where $`p_{}=(p_1,p_2)`$ and the index $`n`$ enumerates the species of the creation and annihilation operators. According to the spectral condition, the operator $`P_{}`$ is positive definite; therefore, the integration over $`p_{}`$ in the given formula is performed only from $`0`$ to $`\mathrm{}`$. The operator $`P_{}`$ vanishes on the physical vacuum $`\mathrm{\Omega }`$, i.e., $`P_{}|\mathrm{\Omega }=0`$, whence it follows that $`a_n(p_{},p_{})|\mathrm{\Omega }=0`$.
For this reason, the physical vacuum $`|\mathrm{\Omega }`$ can be taken as the mathematical vacuum of the Fock space generated by the operators $`a_{n}^{}{}_{}{}^{+}`$. No question of describing the physical vacuum structure arises. The spectrum of the bound states can be found by solving the Schrodinger equation
$$\begin{array}{c}\text{ }P_+|\mathrm{\Psi }=p_+^{}|\mathrm{\Psi }\text{ }\hfill \end{array}$$
under the conditions
$$\begin{array}{c}\text{ }P_{}|\mathrm{\Psi }=p_{}^{}|\mathrm{\Psi },P_{}|\mathrm{\Psi }=0,\text{ }\hfill \end{array}$$
where $`p_+^{}`$, $`p_{}^{}`$ are numbers and the mass squared is defined by
$$\begin{array}{c}\text{ }m^2=2p_+^{}p_{}^{}.\text{ }\hfill \end{array}$$
Here, $`|\mathrm{\Psi }`$ is a vector in the just mentioned Fock space. It is easy to satisfy the conditions $`P_{}|\mathrm{\Psi }=p_{}^{}|\mathrm{\Psi }`$ and $`P_{}|\mathrm{\Psi }=0`$, where $`p_{}^{}`$ is chosen arbitrarily. The problem is to solve the Schrodinger equation. Such an approach is usually called the ”light-front Hamiltonian” approach. Obviously, it can also be used to calculate the scattering matrix.
The scheme described faces considerable problems due to divergences that arise and the lack of the explicit Lorentz invariance. But the possibility of avoiding the direct description of the physical vacuum structure is such a considerable advantage that this approach continues to attract attention. Interest in it has increased with the advent of QCD. This paper is devoted to developing this approach.
We note that the light-front coordinates and the similar ”infinitely large-momentum system” are used in quantum field theory not only in the framework of the Hamiltonian approach based on directly solving the Schrodinger equation. Many results have been obtained by studying the limiting case of fast-moving reference frames in the framework of the explicitly Lorentz-invariant. theory of the scattering matrix or of the Green’s functions . But we consider only the Hamiltonian approach in this paper.
For the theory with the given Lorentz-invariant initial action, constructing the light-front Hamiltonian $`P_+`$ proves a difficult problem. The primary reason is the specific divergences at zero values of tin’ momentum $`p_{}`$ of virtual quanta. In particular, the invariant volume element on the hyperboloid $`p_\mu p^\mu =m^2`$ has the form $`d^2p_{}dp_{}/p_{}`$ and contains $`p_{}`$ in the denominator. The situation becomes more complicated in gauge theories. Even in the first papers on this problem , it became clear that canonical quantization on the surface $`x^+=const`$ could be performed only in the gauge $`A_{}=0`$ or in similar gauges (for example, $`_{}A_{}=0`$). The reason is that second-class constraints arise in the theory and solving them, in particular, requires inverting the covariant derivative $`D_{}=_{}+igA_{}`$. But in the gauge $`A_{}=0`$, the Feynman propagator has an extra term $`p_{}`$ in the denominator, which strengthens the singularities at $`p_{}=0`$ (at least, in the perturbation theory).
As a consequence, a special regularization is required, which consists in ”cutting out” the neighborhood of $`p_{}=0`$ in the momentum space one way or another and which results in breaking the Lorentz invariance (until the regularization is removed). This is unavoidable in the liglit-front Hamiltonian approach. In principle, we can preserve the gauge invariance if, instead of ”cutting out” the neighborhood of $`p_{}=0`$, we restrict the space-time with respect to the coordinate $`x^{}`$ ($`Lx^{}L`$) and impose periodic boundary conditions in $`x^{}`$ on all fields . In this case, the spectrum of the momentum $`P_{}`$ becomes discrete, and the ”zero modes” of the fields $`A`$, i.e., the Fourier modes corresponding to $`p_{}=0`$, are explicitly singled out. To preserve the gauge invariance, these zero modes must be taken into account. It was shown in that secondary second-class constraints arise in the theory, from which we must find the above-mentioned zero modes as functions of the other modes and substitute them in the Hamiltonian. These constraints arc so complicated that solving them proves impossible. Therefore, from the very beginning, we must discard the zero modes in the Lagrangian, which results in breaking the gauge invariance. Therefore, the breakdown of the gauge invariance by a regularization is unavoidable in any case. In what follows, to regularize the theory, we ”cut out” the neighborhood of $`p_{}=0`$. In addition, the conventional UV regularization is, of course, necessary.
It follows from the above discussion that the formal canonical quantization on the hypersurface $`x^+=const`$ can produce a Hamiltonian corresponding to a theory that is not equivalent to the initial Lorentz-invariant one, even in the limit of removing the regularization. As a rule, the equivalence can be provided only by adding nonconventional counterterms to the light-front Hamiltonian.
In the last few years, a number of papers have appeared in which the approximate regularized light-front QCD Hamiltonian was directly constructed . Using the renormalization group theory, the form of the Hamiltonian was adjusted based on the requirement that the result be weakly dependent on the UV cutoff, with the relation to the conventional Lorentz-invariant theory not being traced in detail. Methods for simplifying the Hamiltonian and solving the Schrodinger equation were also proposed. The technique described in these papers is of considerable interest and continues to develop. But it remains unclear to what extent the light-front theory thus obtained corresponds to the conventional Lorentz-invariant QCD.
In this connection, we meet the problem of constructing the light-front Hamiltonian such that it produces a theory equivalent to the conventional Lorentz-invariant one in the limit of removing the regularization. As a necessary condition, we should first provide this equivalence in the framework of the perturbation theory. Then, in particular, we can use the technique described in to simplify the obtained Hamiltonian and the subsequent nonperturbative calculations based on the Schrodinger equation. In doing so, we see the relation to the conventional Lorentz-invariant theory.
The problem of constructing the counterterms for the light-front Hamiltonian, which provide the equivalence of this approach to the Lorentz-invariant one, was investigated in in the framework of the perturbation theory in the coupling constant. In that paper, the authors proposed a method for comparing two perturbation series for the Green’s functions constructed using the light-front Hamiltonian for one and the conventional Lorentz-invariant approach for the other. In addition to the UV regularization, they regularized the singularities at $`p_{}=0`$ using the cutoff $`|p_{}|\epsilon >0`$, i.e., eliminating the Fourier components with $`|p_{}|<\epsilon `$ for every field in the theory. It was revealed that for the required equivalence for the nongauge field theories (in particular, for the Yukawa model), it was only necessary to add a few counterterms to the canonical light-front Hamiltonian. But for gauge theories (both Abelian and non-Abelian) under the given regularization, it proved necessary to introduce an infinite number of counterterms in the Hamiltonian. This situation is closely related to the gauge condition $`A_{}=0`$, which is unavoidable in the canonical light-front quantization. It is known that to correctly construct the perturbation theory in this gauge, it is necessary to use the gauge field propagator in the form proposed by Mandelstam and Leibbrandt
$$\begin{array}{c}\text{ }\frac{i\delta ^{ab}}{k^2+i0}\left(g_{\mu \nu }\frac{k_\mu n_\nu +k_\nu n_\mu }{2k_+k_{}+i0}2\frac{k_+}{n_+}\right)=\frac{i\delta ^{ab}}{k^2+i0}\left(g_{\mu \nu }\frac{k_\mu n_\nu +k_\nu n_\mu }{2(kn^{})(kn)/(nn^{})+i0}2\frac{(kn^{})}{(nn^{})}\right)(1)\text{ }\hfill \end{array}$$
where $`\mu ,\nu =+,,1,2`$; $`k_\pm =(k_0\pm k_3)/\sqrt{2}`$; $`n_{,1,2}=0`$, $`n_{+,1,2}^{}=0`$. The additional pole in $`k_{}`$ in this expression is cut out with the regularization $`|k_{}|\epsilon `$. The distortion arising in this case does not disappear in the limit $`\epsilon 0`$. An infinite number of counterterms is required to compensate this distortion.
The aim of this paper is to overcome this difficulty and obtain the required Hamiltonian with a finite number of counterterms. In such a case, we must change the regularization. For this purpose, we propose shifting the pole with respect to $`k_{}`$ in expression (1. Introduction) from the point $`k_{}=0`$ by changing the Lagrangian such that a small regularizing mass parameter $`\mu ^2`$ is added to the quantity $`2k_+k_{}`$. In this case, the distortions caused by cutting out the interval $`|k_{}|<\epsilon `$ are not so large as in the preceding case, and a finite number of counterterms is sufficient to obtain the correct Hamiltonian in the limit of removing the regularization ($`\epsilon 0`$ and then $`\mu 0`$ together with $`\mathrm{\Lambda }\mathrm{}`$, where $`\mathrm{\Lambda }`$ is a UV regularization parameter). We choose the UV regularization such that, after removing the IR regularization (i.e., after setting $`\epsilon =0`$ and $`\mu =0`$) at the intermediate stages, we obtain the Lorentz-invariant Lagrangian regularized in the UV region. This increases the number of ”ghosts,” but the number of necessary counterterms would otherwise increase sharply. We note that in the gauge $`A_{}=0`$, the one-particle irreducible vertex parts with the upper index ”$``$” do not contribute to the Green’s functions. In what follows, we prove the coincidence of every order of the perturbation theory in the limit of removing the regularization only for the Green’s functions and not for the vertex parts with the index ”$``$”. This is sufficient because the masses of the bound states are determined by the poles of the Green’s functions.
Under the given regularization, the Hamiltonian acts on the space with an indefinite metric, which prevents using the conventional variational principle to solve the Schrodinger equation under the conditions of preserving the regularization. Nevertheless, there exist different variational methods that allow solving this equation.
The chosen regularization breaks the local gauge invariance but preserves the global $`SU(3)`$ invariance. For this reason, the number of necessary counterterms, being finite, is essentially larger than the conventional one. We show that there is a way to choose the counterterms such that the obtained light-front theory is perturbatively equivalent to the initial Lorentz-invariant one in the limit of removing the regularization.
We achieve this goal by starting with the Lagrangian
$$\begin{array}{c}\text{ }L=L_0+L_I(2)\text{ }\hfill \end{array}$$
$$\begin{array}{c}\text{ }L_0=\frac{1}{4}\underset{j=0,1}{}(1)^jf_j^{a,\mu \nu }\left(1+\frac{_{}^2}{\mathrm{\Lambda }_j^2}\frac{_{}^2}{\mathrm{\Lambda }^2}\right)f_{j,\mu \nu }^a+\text{ }\hfill \\ \text{ }+\underset{l=0}{\overset{3}{}}\frac{1}{v_l}\overline{\psi }_l\left(i\gamma ^\mu _\mu M_l\right)\psi _l(3)\text{ }\hfill \end{array}$$
$$\begin{array}{c}\text{ }L_I=c_0_\mu A_\nu ^a^\mu A^{a,\nu }+c_{01}_\mu A_\nu ^aN^{\mu \alpha }_\alpha A^{a,\nu }+c_1_\mu A^{a,\mu }_\nu A^{a,\nu }+c_{11}N^{\mu \alpha }_\mu A_\alpha ^a_\nu A^{a,\nu }+\text{ }\hfill \\ \text{ }+c_{12}N^{\mu \alpha }_\mu A_\alpha ^aN^{\nu \beta }_\nu A_\beta ^a+c_2A_\mu ^aA^{a,\mu }+c_3f^{abc}A_\mu ^aA_\nu ^b^\mu A^{c,\nu }+c_{31}f^{abc}A_\mu ^aA_\nu ^bN^{\alpha \mu }_\alpha A^{c,\nu }+\text{ }\hfill \\ \text{ }+A_\mu ^aA_\nu ^bA_\gamma ^cA_\delta ^d\left(c_4f^{abe}f^{cde}g^{\mu \gamma }g^{\nu \delta }+\delta ^{ab}\delta ^{cd}\left(c_5g^{\mu \gamma }g^{\nu \delta }+c_6g^{\mu \nu }g^{\gamma \delta }\right)\right)+\text{ }\hfill \\ \text{ }+c_7\overline{\psi }\gamma ^\mu i_\mu \psi +c_{71}\overline{\psi }\gamma _\mu N^{\mu \nu }i_\nu \psi +c_{72}\overline{\psi }\gamma _\mu N^{\nu \mu }i_\nu \psi c_8\overline{\psi }\psi +\text{ }\hfill \\ \text{ }+c_9A_\mu ^a\overline{\psi }\gamma ^\mu \frac{\lambda ^a}{2}\psi +c_{91}A_\alpha ^aN^{\mu \alpha }\overline{\psi }\gamma _\mu \frac{\lambda ^a}{2}\psi ,(4)\text{ }\hfill \end{array}$$
where
$$\begin{array}{c}\text{ }f_{j,\mu \nu }^a=_\mu A_{j,\nu }^a_\nu A_{j,\mu }^a,\frac{1}{\mathrm{\Lambda }_j^2}=\{\begin{array}{cc}1/\mathrm{\Lambda }^2,\hfill & j=0\text{,}\hfill \\ 1/\mathrm{\Lambda }^2+1/\mu ^2,\hfill & j=1\text{,}\hfill \end{array}(5)\text{ }\hfill \end{array}$$
$$\begin{array}{c}\text{ }v_0=1,\underset{l=0}{\overset{3}{}}v_l=0,\underset{l=0}{\overset{3}{}}v_lM_l=0,\underset{l=0}{\overset{3}{}}v_lM_l^2=0,(6)\text{ }\hfill \end{array}$$
$$\begin{array}{c}\text{ }A_\mu ^a=\underset{j=0,1}{}A_{j,\mu }^a,\psi =\underset{l=0}{\overset{3}{}}\psi _l.(7)\text{ }\hfill \end{array}$$
Here
$$\begin{array}{c}\text{ }_{}^2=2_+_{},_{}^2=_1^2+_2^2,N^{\alpha \beta }=\frac{n^\alpha n^\beta }{(nn^{})},(8)\text{ }\hfill \end{array}$$
$`\gamma ^\mu `$ are the Dirac matrices, $`\lambda ^a`$ are the matrices of the fundamental representation of the gauge $`SU(3)`$ group, and
$$\begin{array}{c}\text{ }\mathrm{Tr}\lambda ^a=0,\lambda ^a\lambda ^b=if^{abc}\lambda ^c+d^{abc}\lambda ^c+\frac{2}{3}\delta ^{ab},\mathrm{Tr}(\lambda ^a\lambda ^b)=2\delta ^{ab}.(9)\text{ }\hfill \end{array}$$
We assume that the fields $`A_{j,\mu }^a`$, are restricted by the condition
$$\begin{array}{c}\text{ }A_{j,}^a=0.(10)\text{ }\hfill \end{array}$$
In addition, we introduce the cutoff in the momenta $`k_{}`$ and $`k_{}`$
$$\begin{array}{c}\text{ }\epsilon |k_{}|V,v^2k_{}^2V^2(11)\text{ }\hfill \end{array}$$
for all fields from the very beginning. This cutoff as well as condition (1. Introduction) excludes certain degrees of freedom directly in the Lagrangian; such a procedure does not lead to new constraints in the canonical formalism, because no variation with respect to the excluded degrees of freedom is performed.
The quantities $`\epsilon `$, $`\mathrm{\Lambda }`$, $`\mu `$, $`v`$, $`V`$ and $`M_1`$-$`M_3`$ are regularization parameters. The coefficients $`c_i`$ are renormalization constants, i.e., they are functions of the regularization parameters and are expansions in the coupling constant $`g`$. These expansions begin with $`g`$ (with the coefficient one) for $`c_3`$ and $`c_9`$ and with $`g^2`$ and higher for the others. Expression (1. Introduction) contains the conventional QCD interaction and the additional terms necessary for the renormalizability under the assumption that the Lorentz invariance and the global, but not local, gauge invariance are preserved.
In Sec. 2. The light-front Hamiltonian, we prove that the light-front Hamiltonian corresponds to the given Lagrangian. In the framework of the perturbation theory with respect to the coupling constant $`g`$ and with a certain dependence of the renormalization constants on the regularization parameters, this Hamiltonian produces the Green’s functions of the fields $`A_{0,\mu }^a`$, $`\overline{\psi }_0`$, and $`\psi _0`$ coinciding with the Green’s functions of the conventional QCD (renormalized in the Lorentz coordinates using dimensional regularization) in every order with respect to $`g`$ in the limit of removing the regularization (according to the special prescription).
## 2. The light-front Hamiltonian
We develop the canonical light-front formalism for the theory with Lagrangian (1. Introduction)-(1. Introduction) defined on the fields satisfying the condition $`A_{j,}^a=0`$ and subject to cutting out the momenta according to formula (1. Introduction). We use the following representation for the bispinors $`\psi _l`$ and the matrices $`\gamma ^\mu `$:
$$\begin{array}{c}\text{ }\psi =\left(\begin{array}{c}\psi _+\\ \psi _{}\end{array}\right),\gamma ^0=\left(\begin{array}{cc}0& i\\ i& 0\end{array}\right),\gamma ^3=\left(\begin{array}{cc}0& i\\ i& 0\end{array}\right),\gamma ^k=\left(\begin{array}{cc}i\sigma _k& 0\\ 0& i\sigma _k\end{array}\right).(12)\text{ }\hfill \end{array}$$
Here and in what follows, the indices $`k,l`$ ranges over $`1,2`$.
Let us rewrite the expression for the $`L`$ in a form more convenient for the canonical light-front formalism development using instead of the $`A_{j,+}^a`$ new variables
$$\begin{array}{c}\text{ }\phi _j^a=_{}A_{j,+}^a_kA_{j,k}^a.(13)\text{ }\hfill \end{array}$$
For example, for the $`L_I`$ we get:
$$\begin{array}{c}\text{ }L_I=c_{01}^{}A_k^a_{}_+A_k^ac_0A_k^a_{}^2A_k^a+c_1\phi ^a\phi ^a+c_{11}\left(\phi ^a+_kA_k^a\right)\phi ^a+\text{ }\hfill \\ \text{ }+c_{12}\left(\phi ^a+_kA_k^a\right)\left(\phi ^a+_lA_l^a\right)c_2A_k^aA_k^a+c_3f^{abc}A_\mu ^aA_\nu ^b^\mu A^{c,\nu }+c_{31}f^{abc}A_\mu ^aA_\nu ^bN^{\alpha \mu }_\alpha A^{c,\nu }+\text{ }\hfill \\ \text{ }+A_\mu ^aA_\nu ^bA_\gamma ^cA_\delta ^d\left(c_4f^{abe}f^{cde}g^{\mu \gamma }g^{\nu \delta }+\delta ^{ab}\delta ^{cd}\left(c_5g^{\mu \gamma }g^{\nu \delta }+c_6g^{\mu \nu }g^{\gamma \delta }\right)\right)+\text{ }\hfill \\ \text{ }+c_{71}^{}i\sqrt{2}\psi _+^+_+\psi _++c_{72}^{}i\sqrt{2}\psi _{}^+_{}\psi _{}+c_7\left(\psi _+^+i\widehat{}_{}\psi _{}+h.c\right)c_8\left(i\psi _{}^+\psi _++h.c\right)+\text{ }\hfill \\ \text{ }+c_{91}^{}\sqrt{2}A_+^a\psi _+^+\frac{\lambda ^a}{2}\psi _++c_9(\psi _+^+\widehat{A}_{}\psi _{}+h.c.),(14)\text{ }\hfill \end{array}$$
where
$$\begin{array}{c}\text{ }\phi ^a=\underset{j=0,1}{}\phi _j^a,\widehat{A}_{}A_k^a\left(\frac{\lambda ^a}{2}\sigma _k\right),\widehat{}_{}\left(\sigma _k_k\right),\text{ }\hfill \\ \text{ }c_{01}^{}=2c_0+c_{01},c_{71}^{}=c_7+c_{71},c_{72}^{}=c_7+c_{72},c_{91}^{}=c_9+c_{91}.(15)\text{ }\hfill \end{array}$$
We transform the part of the Lagrangian that contains the derivatives $`_+`$ of the fermion fields $`\psi `$, presenting it as a diagonalized bilinear form:
$$\begin{array}{c}\text{ }\underset{l=0}{\overset{3}{}}\frac{1}{v_l}\psi _{l+}^+_+\psi _{l+}+c_{71}^{}\psi _+^+_+\psi _+=\underset{l=0}{\overset{3}{}}B_{ll^{}}\psi _{l+}^+_+\psi _{l^{}+}=\underset{l=0}{\overset{3}{}}\frac{1}{w_l}\psi _{l+}^+_+\psi _{l+}^{},(16)\text{ }\hfill \end{array}$$
where
$$\begin{array}{c}\text{ }\psi _{l+}=\underset{l^{}=0}{\overset{3}{}}U_{ll^{}}\psi _{l^{}+}^{},(17)\text{ }\hfill \end{array}$$
$`1/w_l`$ are the eigenvalues of the matrix $`B_{ll^{}}`$, and $`U_{ll^{}}`$ is the matrix of the transformation diagonalizing $`B_{ll^{}}`$.
Consider the part of the Dirac equation that is a result of the variation with respect to $`\psi _l^+`$:
$$\begin{array}{c}\text{ }\frac{1}{v_l}\left(\sqrt{2}_{}\psi _l+\widehat{}_{}\psi _{l+}M_l\psi _{l+}\right)+c_{72}^{}\sqrt{2}_{}\psi _{}+c_7\widehat{}_{}\psi _+c_8\psi _+ic_9\widehat{A}_{}\psi _+=0.(18)\text{ }\hfill \end{array}$$
This equation does not contain derivatives with respect to $`x^+`$. It is a constraint that can be used to express $`\psi _l`$ in terms of other variables. To do this we sum the equation (2. The light-front Hamiltonian) over index $`l`$ with the weight $`v_l`$. Using the equations (1. Introduction) we obtain
$$\begin{array}{c}\text{ }\sqrt{2}_{}\psi _{}=\underset{l^{}=0}{\overset{3}{}}\left(\widehat{}_{}M_l^{}\right)\psi _{l^{}+}.(19)\text{ }\hfill \end{array}$$
By the substitution of this expression to the equation (2. The light-front Hamiltonian) one can obtain easily from the latter
$$\begin{array}{c}\text{ }\psi _l=\frac{1}{\sqrt{2}}_{}^1Y_l,(20)\text{ }\hfill \end{array}$$
where
$$\begin{array}{c}\text{ }Y_l=\left(\widehat{}_{}M_l\right)\psi _{l+}+v_l\left[c_{72}^{}\underset{l^{}=0}{\overset{3}{}}\left(\widehat{}_{}M_l^{}\right)\psi _{l^{}+}c_7\widehat{}_{}\psi _++c_8\psi _++ic_9\widehat{A}_{}\psi _+\right].(21)\text{ }\hfill \end{array}$$
Let us remark that such a constraint can be resolved in this way not only in $`A_{}=0`$ gauge but in any gauge, if Pauli-Villars fermions are present (if no Pauli-Villars fermions we need to invert the operator $`D_{}=_{}igA_{}`$, and therefore to use the $`A_{}=0`$ gauge).
In turn, the $`\psi _+`$ is expressed in terms of the $`\psi _{l+}^{}`$ by (1. Introduction) and (2. The light-front Hamiltonian). After substituting the expression (2. The light-front Hamiltonian) into Lagrangian the latter depends only on the variables $`A_{j,\mu }^a`$ and $`\psi _{l+}^{}`$.
We eliminate the derivatives $`_+^2A_{j,k}^a`$ in the Lagrangian by integrating by parts and find the momenta conjugate to the $`A_{j,k}^a`$ (by corresponding variation of the $`L`$ at fixed $`\phi _j^a`$):
$$\begin{array}{c}\text{ }\mathrm{\Pi }_{j,k}^a=\frac{L}{(_+A_{j,k}^a)}=\frac{(1)^j}{\mathrm{\Lambda }_j^2}\left[4_{}^2_+A_{j,k}^a_{}^2_{}A_{j,k}^a+\mathrm{\Lambda }_j^2\left(1\frac{_{}^2}{\mathrm{\Lambda }^2}\right)_{}A_{j,k}^a\right]c_{01}^{}_{}A_k^a.(22)\text{ }\hfill \end{array}$$
If we similarly define the momenta conjugate to $`\phi _j^a`$, we obtain the second-class constraint. Using the Legendre transformation with respect to the variables $`_+A_{j,k}^a`$ and $`\mathrm{\Pi }_{j,k}^a`$, which does not affect the variable $`_+\phi _j^a`$, we pass to the first-order Lagrangian. This Lagrangian depends on the derivatives $`_+\psi _l^{}`$, $`_+\phi _j^a`$, and $`_+A_{j,k}^a`$ only linearly (the dependence on $`_+\psi _l^{}`$ and $`_+\phi _j^a`$ is linear from the beginning). Using the Fourier transformation-type formulas, we then pass to the new variables, playing the role of the creation and annihilation operators in order that the following conditions be satisfied. First, the part of the Lagrangian containing derivatives with respect to $`x^+`$ must have a standard canonical form. This automatically solves the above-mentioned constraint for the variable $`\phi _j^a`$. Second, the positive-definite free-part must be diagonal in the corresponding Fock space. The following changes of variables meet these conditions:
$$\begin{array}{c}\text{ }A_{j,k}^a(x)=\frac{1}{(2\pi )^{3/2}}\underset{\epsilon p_{}V}{}𝑑p_{}\underset{v^2p_{}^2V^2}{}d^2p_{}\underset{r=0,1}{}\frac{a_{jr,k}^a(p)e^{ipx}+h.c.}{\sqrt{2\omega _j}},\text{ }\hfill \\ \text{ }\mathrm{\Pi }_{j,k}^a(x)=\frac{i}{(2\pi )^{3/2}}\underset{\epsilon p_{}V}{}dp_{}\underset{v^2p_{}^2V^2}{}d^2p_{}\underset{r=0,1}{}(1)^r\sqrt{\frac{\omega _j}{2}}[a_{jr,k}(p)e^{ipx}h.c.],\text{ }\hfill \\ \text{ }\phi _j^a(x)=\frac{i\mathrm{\Lambda }_j}{(2\pi )^{3/2}}\underset{\epsilon p_{}V}{}𝑑p_{}\underset{v^2p_{}^2V^2}{}d^2p_{}\frac{a_j^a(p)e^{ipx}h.c.}{\sqrt{2p_{}}},\text{ }\hfill \\ \text{ }\psi _{l+,s}^i(x)=\frac{2^{1/4}}{(2\pi )^{3/2}}\underset{\epsilon p_{}V}{}𝑑p_{}\underset{v^2p_{}^2V^2}{}d^2p_{}\left(b_{l,s}^i(p)e^{ipx}+d_{l,s}^{i}{}_{}{}^{+}(p)e^{ipx}\right),(23)\text{ }\hfill \end{array}$$
where $`\omega _j=p_{}\left|1j\frac{p_{}^2}{\mu ^2}\right|=p_{}\left(\frac{p_{}^2}{\mu ^2}1\right)^j`$, $`O(p)O(p_{},p_{})`$, $`pxp_{}x^{}+p_kx^k`$, the index $`s`$ enumerates the spinor components ($`s=1,2`$), and $`i`$ is the index of the fundamental representation of the color group. All creation and annihilation operators are defined for $`p_{}\epsilon `$. We also assume that $`v>\mu `$ in these formulas. The resulting Lagrangian has the form
$$\begin{array}{c}\text{ }L=i\underset{\epsilon p_{}V}{}dp_{}\underset{v^2p_{}^2V^2}{}d^2p_{}\{\underset{j,a}{}(1)^ja_{j}^{a}{}_{}{}^{+}(p)_+a_j^a(p)+\text{ }\hfill \\ \text{ }+\underset{j,r,a,k}{}(1)^ra_{jr,k}^{a}{}_{}{}^{+}(p)_+a_{jr,k}^a(p)+\underset{l,i}{}\frac{1}{w_l}\left(b_{l}^{i}{}_{}{}^{+}(p)_+b_l^i(p)+d_{l}^{i}{}_{}{}^{+}(p)_+d_l^i(p)\right)\}H,(24)\text{ }\hfill \end{array}$$
where $`H=P_+`$ is a Hamiltonian.
Accordingly, the (anti)commutation relations have the form
$$\begin{array}{c}\text{ }[a_{jr,k}^a(p),a_{j^{}r^{},k}^a^{}(p^{})]=[a_{jr,k}^a(p),a_j^{}^a^{}(p^{})]=[a_{jr,k}^a(p),a_{j^{}}^{a^{}}{}_{}{}^{+}(p^{})]=0,\text{ }\hfill \\ \text{ }[a_{jr,k}^a(p),a_{j^{}r^{},k^{}}^{a}{}_{}{}^{+}(p^{})]=(1)^r\delta ^{aa^{}}\delta _{jj^{}}\delta _{rr^{}}\delta _{kk^{}}\delta (p_{}p_{}^{})\delta ^{(2)}(p_{}p_{}^{}),\text{ }\hfill \\ \text{ }[a_j^a(p),a_{j^{}}^{a}{}_{}{}^{+}(p^{})]=(1)^j\delta ^{aa^{}}\delta _{jj^{}}\delta (p_{}p_{}^{})\delta ^{(2)}(p_{}p_{}^{}),\text{ }\hfill \\ \text{ }\{b_{l,s}^i(p),b_{l^{},s^{}}^{i^{}}{}_{}{}^{+}(p^{})\}=\{d_{l,s}^i(p),d_{l^{},s^{}}^{i^{}}{}_{}{}^{+}(p^{})\}=\text{ }\hfill \\ \text{ }=w_l\delta ^{ii^{}}\delta _{ll^{}}\delta _{ss^{}}\delta (p_{}p_{}^{})\delta ^{(2)}(p_{}p_{}^{}).(25)\text{ }\hfill \end{array}$$
The negative signs in the right-hand side of the (anti)commutation relations correspond to the degrees of freedom carrying an indefinite metric in the space of states, and the corresponding operators are ”ghosts”.
In contrast to the conventional canonical light-front formalism, this formalism contains no constraints that halve the number of the Fourier components of the fields $`A_{j,k}^a`$. To preserve the positivity of the momentum $`p_{}`$ everywhere, we are forced to double the number of the creation and annihilation operators $`a_{jr,k}^{a}{}_{}{}^{+}`$ and $`a_{jr,k}^a`$ by introducing the index $`r=0,1`$. The first-order part of the Lagrangian, which contains no derivatives with respect to $`x^+`$, coincides with the Hamiltonian $`H=P_+`$ up to a sign. It has the form
$$\begin{array}{c}\text{ }H=𝑑x^{}d^2x^{}\underset{j=0,1}{}\{\frac{(1)^j\mathrm{\Lambda }_j^2}{8}\left[(1)^j\left(_{}^1\mathrm{\Pi }_{j,k}^a+c_{01}^{}A_k^a\right)+\frac{_{}^2}{\mathrm{\Lambda }_j^2}A_{j,k}^a\left(1\frac{_{}^2}{\mathrm{\Lambda }^2}\right)A_{j,k}^a\right]^2+\text{ }\hfill \\ \text{ }+\frac{(1)^{j+1}}{2}A_{j,k}^a_{}^2\left(1\frac{_{}^2}{\mathrm{\Lambda }^2}\right)A_{j,k}^a+\frac{(1)^{j+1}}{2}\phi _j^a\left(1\frac{_{}^2}{\mathrm{\Lambda }^2}\right)\phi _j^a\}+c_0A_k_{}^2A_k+\text{ }\hfill \\ \text{ }+\frac{i}{\sqrt{2}}\underset{l=0}{\overset{3}{}}\frac{1}{v_l}\psi _{l+}^+\left(_{}^2M_l^2\right)_{}^1\psi _{l+}\frac{i}{\sqrt{2}}c_{72}^{}\underset{l,l^{}=0}{\overset{3}{}}\psi _{l+}^+\left(\widehat{}_{}+M_l\right)_{}^1\left(\widehat{}_{}M_l^{}\right)\psi _{l^{}+}+(26)\text{ }\hfill \\ \text{ }+(\frac{i}{\sqrt{2}}c_7\psi _+^+\widehat{}_{}_{}^1\underset{l=0}{\overset{3}{}}\left(\widehat{}_{}M_l\right)\psi _{l+}+h.c.)+(\frac{i}{\sqrt{2}}c_8\psi _+^+_{}^1\underset{l=0}{\overset{3}{}}\left(\widehat{}_{}M_l\right)\psi _{l+}+h.c.)+\text{ }\hfill \\ \text{ }+\sqrt{2}c_{91}^{}\left(\phi ^a+_kA_k^a\right)_{}^1\left(\psi _+^+\frac{\lambda ^a}{2}\psi _+\right)+(\frac{1}{\sqrt{2}}c_9\psi _+^+\widehat{A}_{}_{}^1\underset{l=0}{\overset{3}{}}\left(\widehat{}_{}M_l\right)\psi _{l+}+h.c.)L_I^{}.\text{ }\hfill \\ \text{ }\hfill \end{array}$$
Here, $`L_I^{}`$ denotes expression (2. The light-front Hamiltonian) with the terms with the coefficients $`c_0`$, $`c_{01}^{}`$, $`c_7`$, $`c_{71}^{}`$, $`c_{72}^{}`$, $`c_8`$, $`c_9`$, $`c_{91}^{}`$ omitted, $`\psi _{l+}`$ is expressed in terms of $`\psi _{l+}^{}`$ by formulas (1. Introduction) and (2. The light-front Hamiltonian), and the quantities $`A_{j,k}^a`$, $`\mathrm{\Pi }_{j,k}^a`$, $`\phi _j^a`$, $`\psi _{l+}^{}`$ and $`\psi _{l+}^+`$ are expressed in terms of the creation and annihilation operators by formulas (2. The light-front Hamiltonian).
The operator $`P_{}`$ has the form
$$\begin{array}{c}\text{ }P_{}=\underset{\epsilon p_{}V}{}dp_{}\underset{v^2p_{}^2V^2}{}d^2p_{}p_{}\{\underset{j=0,1}{}(1)^j[a_j^a{}_{}{}^{+}(p)a_j^a(p)+\text{ }\hfill \\ \text{ }+\underset{r=0,1}{}(1)^{j+r}a_{jr,k}^{a}{}_{}{}^{+}(p)a_{jr,k}^a(p)]+\text{ }\hfill \\ \text{ }+\underset{l=0}{\overset{3}{}}\frac{1}{w_l}\left(b_{l}^{i}{}_{}{}^{+}(p)b_l^i(p)+d_{l}^{i}{}_{}{}^{+}(p)d_l^i(p)\right)\}.(27)\text{ }\hfill \end{array}$$
This operator is positive definite in the Fock space with the vacuum defined by
$$\begin{array}{c}\text{ }a_j^a|0=a_{jr,k}^a|0=b_l^i|0=d_l^i|00.(28)\text{ }\hfill \end{array}$$
In the framework of the perturbation theory, in the limit of removing the regularization, all ”ghosts” are switched off in the sense that the Green’s functions of this theory tend to the Green’s functions of the correct theory, which has no ”ghosts”. This gives us hope that in the limit of removing the regularization, the unitarity condition also holds beyond the scope of the perturbation theory.
## 3. Comparison of the light-front and Lorentz-invariant <br>perturbation theories
By the result of the perturbation theory, we mean the set of the Green’s functions for the fields $`A_{0,\mu }^a`$, $`\overline{\psi }_0`$, and $`\psi _0`$ constructed perturbatively in the coupling constant $`g`$. We regard the fields $`A_{1,\mu }^a`$, $`\overline{\psi }_l`$ and $`\psi _l`$ for $`l=1,2,3`$ as auxiliary. We show that Hamiltonian (2. The light-front Hamiltonian) with a certain $`\mathrm{\Lambda }`$ dependence of the coefficients $`c_i`$ produces a perturbation theory coinciding in the limit of removing the regularization with the renormalized perturbation theory obtained from the conventional QCD Lagrangian in the gauge $`A_{}=0`$ using the Mandelstam-Leibbrandt prescription and dimensional regularization (the conventional perturbation theory in what follows). (See for the renormalization of the conventional perturbation theory.) The regularization is removed as follows: first, $`V\mathrm{}`$, then $`\epsilon 0`$, then $`\mathrm{\Lambda }\mathrm{}`$; $`M_1`$, $`M_2`$, $`M_3`$, $`\mu `$ and $`v`$ are assumed to be functions of $`\mathrm{\Lambda }`$ such that $`M_1,M_2,M_3\mathrm{}`$, $`\mu 0`$ and $`v0`$ as $`\mathrm{\Lambda }\mathrm{}`$. The latter two functions must satisfy more exact restrictions: $`v>\mu `$ (this condition was used in constructing the Hamiltonian) and $`\mu \mathrm{\Lambda },v\mathrm{\Lambda }\underset{\mathrm{\Lambda }\mathrm{}}{}0`$ (these restrictions are obtained below).
First, the noncovariant perturbation theory produced by the Hamiltonian can be obtained from the Feynman perturbation theory constructed based on the Lagrangian corresponding to the given Hamiltonian by resumming diagrams in every order and using the following integration rule: in calculating the diagrams, we first integrate over $`k_+`$ (the momentum component conjugate to the light-front time $`x^+`$) and then over the other components . Therefore, it is sufficient to prove that the perturbation theory obtained from Lagrangian (1. Introduction)-(1. Introduction) and supplemented by the given integration rule and the limiting transitions coincides with the conventional perturbation theory.
We represent the free part of the Lagrangian, Eq. (1. Introduction), in a form convenient for the perturbation theory analysis,
$$\begin{array}{c}\text{ }L_0=\frac{1}{2}\underset{j=0,1}{}(1)^jA_{j,\mu }^a\left(1+\frac{_{}^2}{\mathrm{\Lambda }_j^2}\frac{_{}^2}{\mathrm{\Lambda }^2}\right)\left(g^{\mu \nu }^2^\mu ^\nu \right)A_{j,\nu }^a+\text{ }\hfill \\ \text{ }+\underset{l=0}{\overset{3}{}}\frac{1}{v_l}\overline{\psi }_l\left(i\gamma ^\mu _\mu M_l\right)\psi _l,(29)\text{ }\hfill \end{array}$$
and we separate the part corresponding to the conventional interaction from (1. Introduction),
$$\begin{array}{c}\text{ }L_I^{conv}=gf^{abc}A_\mu ^aA_\nu ^b^\mu A^{c,\nu }\frac{g^2}{4}f^{abe}f^{cde}A_\mu ^aA_\nu ^bA^{c,\mu }A^{d,\nu }gA_\mu ^a\overline{\psi }\gamma ^\mu \frac{\lambda ^a}{2}\psi .(30)\text{ }\hfill \end{array}$$
We assume that the remaining part of (1. Introduction) (where all coefficients are the expansions in $`g`$ starting from the terms of order $`g^2`$) consists of the renormalization counterterms, i.e., it eliminates the divergences arising in the perturbation theory as $`\mathrm{\Lambda }\mathrm{}`$. The notation used and the additional conditions adopted are given by formulas (1. Introduction)-(1. Introduction) and in the text following them.
The propagators of the fields $`A_{j,\mu }^a`$ and $`\psi _l`$ in the momentum space are
$$\begin{array}{c}\text{ }\mathrm{\Delta }_{j,\rho \nu }^{ab}=\frac{i\delta ^{ab}}{k^2+i0}\left(g_{\rho \nu }\frac{k_\rho n_\nu +k_\nu n_\rho }{k_{}}\right)\frac{(1)^j}{1\frac{k_{}^2}{\mathrm{\Lambda }_j^2}+\frac{k_{}^2}{\mathrm{\Lambda }^2}i0},(31)\text{ }\hfill \end{array}$$
$$\begin{array}{c}\text{ }\mathrm{\Delta }_l^\psi =iv_l\frac{\gamma ^\mu k_\mu +M_l}{k^2M_l^2+i0},(32)\text{ }\hfill \end{array}$$
where $`n_+=1`$ and $`n_{},n_{}=0`$.
Because the fields $`A_{j,\mu }^a`$ and $`\psi _l`$ always enter interaction (1. Introduction) in terms of the sums $`A_\mu ^a`$ and $`\psi `$ the sums of the propagators
$$\begin{array}{c}\text{ }\mathrm{\Delta }_{\rho \nu }^{ab}\underset{j}{}\mathrm{\Delta }_{j,\rho \nu }^{ab}\mathrm{and}\mathrm{\Delta }^\psi \underset{l}{}\mathrm{\Delta }_l^\psi \text{ }\hfill \end{array}$$
always enter the diagrams as in the Pauli-Villars regularization. After all diagrams are presented in terms of the summary propagators, we can take the limit $`V\mathrm{}`$ for an arbitrary diagram (i.e., remove the restrictions $`|k_{}|V`$ and $`k_{}^2V^2`$) because, in view of conditions (1. Introduction), the propagator $`\mathrm{\Delta }^\psi `$ decreases sufficiently fast (the sufficiently fast decrease of propagators (3. Comparison of the light-front and Lorentz-invariant perturbation theories) is provided by the finiteness of $`\mathrm{\Lambda }`$) and the integrals converge.
The summary propagator $`\mathrm{\Delta }_{\rho \nu }^{ab}`$ is
$$\begin{array}{c}\text{ }\mathrm{\Delta }_{\rho \nu }^{ab}=\frac{i\delta ^{ab}}{k^2+i0}\left(\frac{k_{}^2}{k_{}^2\widehat{\mu }^2+i0}g_{\rho \nu }\frac{k_\rho n_\nu +k_\nu n_\rho }{k_{}^2\widehat{\mu }^2+i0}2k_+\right)R,(33)\text{ }\hfill \end{array}$$
where
$$\begin{array}{c}\text{ }\widehat{\mu }^2=\mu ^2\frac{\mathrm{\Lambda }^2+k_{}^2}{\mathrm{\Lambda }^2+\mu ^2},R=\frac{1}{\left(1\frac{k^2}{\mathrm{\Lambda }^2}i0\right)\left(1+\frac{\mu ^2}{\mathrm{\Lambda }^2}\right)}.(34)\text{ }\hfill \end{array}$$
After removing the cutoff ($`\mathrm{\Lambda }\mathrm{}`$), we have $`\widehat{\mu }0`$, and the propagator takes the conventional form. In terms of the propagators $`\mathrm{\Delta }_{\rho \nu }^{ab}`$ and $`\mathrm{\Delta }^\psi `$ and the vertices from (3. Comparison of the light-front and Lorentz-invariant perturbation theories), the set of Feynman diagrams is the same as in the conventional perturbation theory.
It is easy to see that as long as $`\mathrm{\Lambda }`$ is finite, there are no UV divergences in the perturbation theory constructed using the Lagrangian under consideration. Using this fact, as well as the condition $`k_{}^2v^2`$ (see (1. Introduction)), we can apply the formalism presented in to the perturbation theory and show that for the majority of diagrams, after the limit $`\epsilon 0`$ is taken, the result of their light-front calculation (where we first integrate over $`k_+`$ and then over the other components according to the rules providing the correspondence with the noncovariant perturbation theory as explained above) coincides with the result of calculating the same diagrams in the Lorentz coordinates. The possible discrepancy that arises in some cases can be compensated by redefining the coefficient $`c_2`$ in the Lagrangian. This is shown in Appendix 1. Therefore, it is sufficient to prove that in the limit $`\mathrm{\Lambda }\mathrm{}`$, the perturbation theory with the summary propagators $`\mathrm{\Delta }_{\rho \nu }^{ab}`$ and $`\mathrm{\Delta }^\psi `$, with interaction (1. Introduction), and with the restrictions $`\epsilon |k_{}|V`$ and $`k_{}^2V^2`$ removed coincides with the conventional perturbation theory. We emphasize that we can now perform all calculations in the Lorentz coordinates; therefore, we can make the Wick rotation and pass to calculating the diagrams in the Euclidean space (the location of the poles of propagators (3. Comparison of the light-front and Lorentz-invariant perturbation theories) allows this).
Propagator (3. Comparison of the light-front and Lorentz-invariant perturbation theories) differs from the propagator of the conventional perturbation theory, first, in the factor $`R`$ providing the UV regularization, second, in the quantity $`\widehat{\mu }`$, which vanishes as $`\mathrm{\Lambda }\mathrm{}`$ and, third, in the condition $`k_{}^2v^2`$, where $`v0`$ as $`\mathrm{\Lambda }\mathrm{}`$.
We now analyze the behavior of an arbitrary Feynman diagram as $`\mu 0`$ and $`v0`$ (for finite $`\mathrm{\Lambda }`$). In this case, essential IR divergences (essential in the sense that they arise for any values of external momenta and not just for special values) can occur. If such a divergence does not appear, then in investigating the limit $`\mathrm{\Lambda }0`$ for an arbitrary diagram, we can at once set $`\widehat{\mu }=0`$ and $`v=0`$ in its integrand. In this case, the error in the integrand contains the factor $`\widehat{\mu }^2`$. The UV divergence of the initial diagram is not worse than quadratic; therefore, after separating the factor $`\widehat{\mu }^2`$, where $`\widehat{\mu }^2\mu ^2(1+k_{}^2/\mathrm{\Lambda }^2)`$, the divergence is not worse than logarithmic, and the condition for an error decrease is $`\mu ^2\mathrm{ln}\mathrm{\Lambda }\underset{\mathrm{\Lambda }\mathrm{}}{}0`$. This consideration does not take an increase in the IR divergence after separating the factor $`\widehat{\mu }^2`$ into account. Its power increases by two. Because there was no divergence before, this power becomes not, greater than one, i.e., integrating the IR divergence gives (in view of the factor $`\widehat{\mu }^2`$) the order $`\widehat{\mu }`$. Integrating over the remaining variables produces an UV divergence not worse than linear. Consequently, the condition for the error decrease is $`\mu \mathrm{\Lambda }\underset{\mathrm{\Lambda }\mathrm{}}{}0`$. Similar considerations give the condition $`v\mathrm{\Lambda }\underset{\mathrm{\Lambda }\mathrm{}}{}0`$.
We analyze the occurrence of essential IR divergences as $`\mu ,v0`$ in Appendix 2. It turns out that such divergences arise only in one case – for the diagram terms which contain the following factor
$$\begin{array}{c}\text{ }\mathrm{\Delta }_{+\rho }G^{\rho \nu }\mathrm{\Delta }_{\nu \alpha },(35)\text{ }\hfill \end{array}$$
where $`G^{\rho \nu }(k)`$ is an arbitrary one-particle irreducible two-point subdiagram, and $`\mathrm{\Delta }`$ are propagators. In which connection the divergence takes place only when index $`\mu `$ in formula (3. Comparison of the light-front and Lorentz-invariant perturbation theories) takes the values $`1,2`$, and the divergence the whole of diagram at that is logarithmic: $`(\mathrm{ln}\mu )^N`$, where $`N`$ is not larger then the number of factors of form (3. Comparison of the light-front and Lorentz-invariant perturbation theories).
It is clear why such a divergence does not produce any problems in the conventional perturbation theory, when $`G^{\rho \nu }(k)=G_{dim}^{\rho \nu }(k)`$ is calculated gauge invariantly with the aid of dimensional regularization. In the expression (3. Comparison of the light-front and Lorentz-invariant perturbation theories) at $`\rho =\mathrm{"}\mathrm{"}`$ the first propagator turns to zero, and, besides, $`g_{++}=g_+=0`$, hence, one should consider in the propagator only the item containing the sum $`(k_+n_\rho +k_\rho n_+)`$. The first term of this sum does not give a divergence because $`n_1=n_2=0`$, and the second term gives the factor $`k_\rho G^{\rho \nu }(k)`$, which is equal to zero because of the Ward identities (analogous to ones adduced in work ) which are the consequences of exact maintenance of gauge invariance. From the given reasoning it is clear that breaking the gauge invariance without the simultaneous regularization of the essential IR divergences makes the perturbation theory senseless. In our case, this regularization is provided by introducing the quantity $`\mu `$.
By induction on the loop number, we prove that with a certain choice of the coefficients of the counterterms in Lagrangian (1. Introduction) in every order, the value of every Feynman diagram tends to its conventional value calculated using dimensional regularization as $`\mathrm{\Lambda }\mathrm{}`$. It is clear that in the one-loop order, there are no subdiagrams; therefore, there are no essential IR divergences. Then on considering a diagram containing factor (3. Comparison of the light-front and Lorentz-invariant perturbation theories), for the lower order $`G^{\rho \nu }(k)`$ subdiagram, we have
$$\begin{array}{c}\text{ }G^{\rho \nu }(k)G_{dim}^{\rho \nu }(k)=O\left(\frac{1}{\mathrm{\Lambda }}\right)(36)\text{ }\hfill \end{array}$$
at $`\rho ,\nu \mathrm{"}\mathrm{"}`$ (just this estimate is obtained in Appendix 4). Because the value $`G^{\rho \nu }(k)`$ at $`\rho =\mathrm{"}\mathrm{"}`$ or $`\nu =\mathrm{"}\mathrm{"}`$ does not give a contribution to the expression (3. Comparison of the light-front and Lorentz-invariant perturbation theories) then taking into account (3. Comparison of the light-front and Lorentz-invariant perturbation theories) we can maintain that accurate to $`O\left(\frac{1}{\mathrm{\Lambda }}\right)`$ the expression (3. Comparison of the light-front and Lorentz-invariant perturbation theories) coincides with
$$\begin{array}{c}\text{ }\mathrm{\Delta }_{+\rho }G_{dim}^{\rho \nu }\mathrm{\Delta }_{\nu \alpha },(37)\text{ }\hfill \end{array}$$
where, as it was already said, the divergence is absent.
Therefore if $`(\mathrm{ln}\mu )^N/\mathrm{\Lambda }\underset{\mathrm{\Lambda }\mathrm{}}{}0`$ then in the limit $`\mathrm{\Lambda }\mathrm{}`$ the diagrams containing the factor of form (3. Comparison of the light-front and Lorentz-invariant perturbation theories) will not differ from their values calculated in the conventional perturbation theory under dimensional regularization. More exactly a possible difference is due to the diagram divergence, but it is polynomial and is compensated by the counterterms of the same form that arise under the renormalization. Therefore, it is now sufficient to prove that in the limit $`\mathrm{\Lambda }\mathrm{}`$, the Euclidean perturbation theory with propagator (3. Comparison of the light-front and Lorentz-invariant perturbation theories), where we set $`\mu =0`$, with the propagator $`\mathrm{\Delta }^\psi `$, with interaction (1. Introduction), and with restrictions (1. Introduction) removed coincides with the conventional perturbation theory.
Propagator (3. Comparison of the light-front and Lorentz-invariant perturbation theories) with $`\mu =0`$ after the transition to the Euclidean space can be written down as
$$\begin{array}{c}\text{ }\mathrm{\Delta }_{\rho \nu }^{ab}=\frac{\delta ^{ab}}{k^2}\left(\delta _{\rho \nu }\frac{k_\rho n_\nu +k_\nu n_\rho }{(kn)}\right)\frac{1}{1+\frac{k^2}{\mathrm{\Lambda }^2}}.(38)\text{ }\hfill \end{array}$$
This expression is Lorentz-invariant if we assume that the vector $`n_\nu `$ (complex in the Euclidean space, such that $`in_0n_3=0`$, $`n_{1,2}=0`$) to be properly transformed under the Lorentz transformation. It is interesting that there is no distinguished vector other than $`n_\nu `$ in the Euclidean space, whereas in the pseudo-Euclidean space, the Mandelstam-Leibbrandt prescription distinguishes the additional surface $`k_+=0`$. However the vector $`n_\nu `$ is complex, and there is new fixed vector, namely, the complex conjugated to $`n_\nu `$. It is seen that it coincides (up to a factor) with Euclidean continuation of the vector $`n_\nu ^{}`$, which is defined below the equation (1. Introduction) and picks out the surface $`k_+=0`$ in pseudoeuclidean space.
Interaction Lagrangian (3. Comparison of the light-front and Lorentz-invariant perturbation theories) is also Lorentz-invariant; therefore, the counterterms that must be added to the Lagrangian under the renormalization in every order of the perturbation theory are Lorentz-invariant. It is shown in the Appendix 4 that the values which it is necessary to add to the diagrams in order to in the limit $`\mathrm{\Lambda }\mathrm{}`$ make them finite and coincident with their values calculated using dimensional regularization are polynomials with respect to external momenta with the coefficients containing factors $`N_{\alpha \beta }`$ (see the definition of this value in formula (1. Introduction)). We should take into account that in the counterterms the vector $`n_\nu `$ cannot be contracted with the field $`A_\nu `$, because we consider the Feynman diagrams for the Green’s functions, whose external lines cannot carry the upper index ”$``$” as it gives zero after convolution with the propagators. It is evident that the counterterms are globally gauge invariant because the initial Lagrangian has this property. We now analyze the possible form of the counterterms.
There exist logarithmically divergent diagrams with four gluon external lines. In Appendix 3, we show that the structure of an arbitrary diagram of this type with respect to the labels of the gauge group can only have the form $`f^{abe}f^{cde}`$ or $`\delta ^{ab}\delta ^{cd}`$. Taking the Lorentz invariance into account, we conclude that the general form of the corresponding counterterms is exhausted by the terms with the coefficients $`c_4`$, $`c_5`$, and $`c_6`$ in Lagrangian (1. Introduction) and we can therefore replace the addition of the counterterms by a certain choice of these coefficients.
There also exist divergent diagrams with three gluon external lines. In general, they can diverge linearly; however, because of the Lorentz invariance, the divergent part must contain the factor $`k^\mu `$, and the divergence is really logarithmic. We show in Appendix 3 that the structure of an arbitrary diagram of this type with respect to the labels of the gauge group can only have the form $`f^{abc}`$. For the general form of the divergence, this gives the terms with the coefficients $`c_3`$ and $`c_{31}`$ in (1. Introduction) (note, that the existence of the latter takes account of the volume $`N_{\alpha \beta }`$ appearance in the counterterms), and we can therefore replace the addition of counterterms by a certain choice of these coefficients. In addition, there exist logarithmically divergent diagrams with two fermion and one gluon external lines. It is evident that the general form of the divergence is defined by the terms with the coefficients $`c_9`$ and $`c_{91}`$ in (1. Introduction) and we can therefore replace the addition of counterterms by a certain choice of these coefficients.
Next, there exist divergent diagrams with two gluon external lines. They diverge quadratically, and the general form of the divergence is given by the term $`c_2`$ in (1. Introduction). But after subtracting the quadratic divergence, a linear divergence can remain. Because of the Lorentz invariance, it is really absent, and only the logarithmic divergence exists. It is evident that the general form of this divergence is given by terms with the coefficients $`c_0`$, $`c_{01}`$, $`c_1`$, $`c_{11}`$, $`c_{12}`$ in (1. Introduction) and we can therefore replace the addition of counterterms by a certain choice of these coefficients and coefficient $`c_2`$. There also exist divergent diagrams with two fermion external lines. They diverge linearly, and the general form of the divergence is given by the term with the coefficient $`c_8`$ in (1. Introduction). But after subtracting the linear divergence, a logarithmic divergence can remain. It is evident that the general form of this divergence is given by the term with the coefficients $`c_7`$, $`c_{71}`$, $`c_{72}`$ in (1. Introduction) and we can therefore replace the addition of counterterms by a choice of these coefficients and coefficient $`c_8`$.
We conclude that the perturbation theory with propagator (3. Comparison of the light-front and Lorentz-invariant perturbation theories), with the propagator $`\mathrm{\Delta }^\psi `$ and with interaction (1. Introduction) (we must compare such a perturbation theory with the conventional one) is renormalizable by renormalizing the coefficients $`c_i`$. It is clear that by properly adjusting the additions to the quantities $`c_i`$ in every order of the perturbation theory, i.e., by manipulating the finite renormalizations, in the limit $`\mathrm{\Lambda }\mathrm{}`$, we can achieve the coincidence of the value of every Feynman diagram with its conventional value (the corresponding scheme is briefly presented in Appendix 4). This is just what we wanted to prove.
Acknowledgments. This work was supported in part (S. A. P.) by the grant 96-0457 INTAS within the framework of the research program of International Center of Fundamental Physics in Moscow (ICFPM) and the Euler Program of Berlin Free University.
## Appendix 1
We compare the results of calculating an arbitrary Feynman diagram in the light-front coordinates (with the limit $`\epsilon 0`$ subsequently taken) and in the Lorentz coordinates. Every diagram is constructed from summary propagators (3. Comparison of the light-front and Lorentz-invariant perturbation theories) and $`\mathrm{\Delta }^\psi `$ as well as from the vertices entering interaction Lagrangian (1. Introduction) with the conditions $`|k_{}|\epsilon `$ and $`k_{}^2v^2`$ but without the conditions $`|k_{}|V`$ and $`k_{}^2V^2`$.
We use the formalism presented in . Under the condition $`k_{}^2v^2`$, which is equivalent to the presence of the nonzero mass in two dimensions, the form of propagator (3. Comparison of the light-front and Lorentz-invariant perturbation theories) is admissible for this formalism. It is easy to see that for all diagrams, the index $`\omega _{}`$ of the UV divergence with respect to $`k_+`$ and $`k_{}`$ is negative; therefore, for our theory, there are no special cases described in . The numerators of all integrands of the Feynman diagrams are polynomials; therefore, we have $`\tau >0`$ and $`\eta >0`$ (for the notation, see ). The basic formula is
$$\begin{array}{c}\text{ }\sigma =min(\tau ,\omega _{}\omega _+\mu +\eta ),(\mathrm{A1}.1)\text{ }\hfill \end{array}$$
where the minimum is taken over all subdiagrams, and the required difference between the light-front and Lorentz-invariant calculations is of the order $`\epsilon ^\sigma `$.
An external gluon line carrying the label ”$`+`$” contributes $`+1`$ to the value of $`\omega _{}\omega _+`$. A pair of external fermion lines that are connected with the continuous fermion line contributes $`1`$ (if the diagram is proportional to $`\gamma ^+`$ with respect to the labels of this pair) or $`+1`$ (if it, is proportional to $`\gamma ^{}`$) or $`0`$ (in the other cases) to the value of $`\omega _{}\omega _+`$. We can see from formula (3. Comparison of the light-front and Lorentz-invariant perturbation theories) that without considering the factors from the vertices, the summary gluon propagator of the $`\mathrm{\Pi }`$-line contributes $`2`$ (if the propagator carries the labels ”$``$”) or $`3`$ (if the propagator carries the labels ”$`+`$” or ”$`++`$”) to the value of $`\mu `$. The factors from the vertices carrying the index ”$`+`$” contribute $`+1`$ to the value of $`\mu `$. From formula (3. Comparison of the light-front and Lorentz-invariant perturbation theories) with conditions (1. Introduction) taken into account, we see that the summary fermion propagator of the $`\mathrm{\Pi }`$-line contributes $`2`$ to the value of $`\mu `$. This number increases if we decrease the number of the additional Pauli-Villars fermion fields. Analyzing this information, we obtain the general form of the diagrams with $`\sigma 0`$ (in fact, $`\sigma =0`$). It is presented in Fig. 1.
The conditions for these diagrams to be trivially dependent on external momenta are fulfilled (see ); therefore, for these diagrams, the required difference is
$$\begin{array}{c}\text{ }Cg_{AB},(\mathrm{A1}.2)\text{ }\hfill \end{array}$$
where $`C`$ is a constant in the limit $`\epsilon 0`$, i.e., the dependence on external momenta is absent. There is no possible logarithmic correction $`\mathrm{ln}(\epsilon /u)`$, because, in view of the Lorentz invariance in the space of $`k_+`$ and $`k_{}`$, which holds for the Lagrangian, the quantity $`u`$ must behave like the ”$``$”-component of a vector, but there are no such expressions. There could only be the component of an external momentum, but, we already know that there is no dependence on it. Therefore, all the difference is compensated by adding the term of the form
$$\begin{array}{c}\text{ }A_\mu ^ag^{\mu \nu }A_\nu ^a,(\mathrm{A1}.3)\text{ }\hfill \end{array}$$
to the Lagrangian, i.e., by redefining the coefficient $`c_2`$ in formula (1. Introduction).
## Appendix 2
We analyze the possibility of the occurrence of the essential IR divergences (i.e., those occurring at any value of the external momenta) in the Feynman diagrams in the Euclidean space. We consider two cases of the divergences: when integrating over only the longitudinal momenta (they can appear because of (the term proportional to $`k_+/k_{}^2`$ in the gluon propagator) and when integrating over all momenta (the factors $`1/k^2`$ in the propagators also contribute to these divergences).
First case. We study whether a divergence exists if a part (or all) of the longitudinal loop momenta tend to some finite values. We assume that the external momenta do not take the special values (where the sum of a part of them is equal to zero). In addition, we assume that the transverse loop momenta are such that for all of the propagator momenta, we have $`Q_{}0`$ (if this condition is violated, we must take the extra contribution from $`d^2Q_{}`$ into account; this is the second case). Then the divergence can arise only if for some gluon line, we have $`Q_{}=0`$. The factor in the integrand producing the possible divergence has the form
$$\begin{array}{c}\text{ }\frac{Q_\mu n_\nu +Q_\nu n_\mu }{Q_\alpha n^\alpha }=\sqrt{2}\frac{(Q_\mu n_\nu +Q_\nu n_\mu )(iQ_0+Q_3)}{Q_0^2+Q_3^2}(\mathrm{A2}.1)\text{ }\hfill \end{array}$$
and is a pole of an order not higher than one. In the loop momentum space, we consider a point where $`Q_{}=0`$ for a certain set of lines. We look for the power $`\sigma `$ of the IR divergence. Every line of this set contributes $`1`$ (if it carries the labels $`+`$, see (Appendix 2)) or $`0`$ (if it carries the labels $``$ or $`++`$). We exclude all lines with the labels $``$ and $`++`$; then every line of the set contributes $`1`$. The differentials $`d^2q_{}`$ of the loop momenta (the integration volume elements) give a positive contribution to $`\sigma `$. We must consider only those loop momenta whose change (other momenta being fixed) violates the condition $`Q_{}=0`$ for the lines of the set. The number of such loop momenta is the number of lines whose momenta can be taken arbitrarily, i.e., the number of independent lines. The total positive contribution to $`\sigma `$ is equal to twice this number. We then find other contributions to the IR divergence.
We break all lines of the set. The diagram splits into $`n+1`$ connected parts ($`n=0,1,\mathrm{}`$). All momenta external with respect to the whole diagram enter one part (the external momenta would otherwise take the special values). We call this part separated and the other parts nonseparated.
A nonseparated part is a subdiagram; if it has the external Lorentz labels, then it must be proportional to its external line momenta carrying these labels. But the factor of proportionality cannot contribute to $`\sigma `$, because of the Lorentz invariance, the invariance with respect to multiplying $`n^\mu `$ by a complex number, and the condition $`Q_{}0`$. Every external line carrying the label $`k`$ ($`k=1,2`$) gives the factor $`k^k`$ or $`g^{\mu k}`$; the line carrying the label ”$`+`$” gives the factor $`k^+`$, where $`k`$ is a linear combination of the external momenta of the subdiagram, i.e., the momenta of the set lines. Therefore, every external line carrying the label ”$`+`$” contributes $`+1`$ to $`\sigma `$, and one carrying the label $`k`$ contributes zero. Let $`m`$ be a summary contribution to $`\sigma `$ obtained in such a way from the nonseparated parts.
Let $`s`$ be the number of lines of the set and $`r`$ be the number of lines of the set external with respect to the separated part. Every line carries the label ”$`+`$” at one end; therefore, we have
$$\begin{array}{c}\text{ }msr.(\mathrm{A2}.2)\text{ }\hfill \end{array}$$
The number of the independent lines in the set is equal to $`sn`$ (each of the $`n+1`$ parts gives the $`\delta `$-function, and one $`\delta `$-function is common to the whole diagram). By the definition of $`\sigma `$, we have
$$\begin{array}{c}\text{ }\sigma =2(sn)+ms.(\mathrm{A2}.3)\text{ }\hfill \end{array}$$
Using (Appendix 2), we obtain
$$\begin{array}{c}\text{ }\sigma 2(sn)r.(\mathrm{A2}.4)\text{ }\hfill \end{array}$$
Every nonseparated part has a minimum of two external lines. This means that all parts (separated and nonseparated) together have a minimum of $`2n+r`$ external lines of the set. All these lines are pairwise connected with each other. Consequently, we have
$$\begin{array}{c}\text{ }s\frac{1}{2}(2n+r)=n+\frac{r}{2},(\mathrm{A2}.5)\text{ }\hfill \end{array}$$
whence we obtain $`\sigma 0`$, i.e., the divergence cannot be worse than logarithmic.
We find the general case producing this divergence. In this case, all the above-cited inequalities must reduce to equalities, i.e., we have $`m=sr`$, and all the nonseparated parts are two-point diagrams. The general form of such a diagram is shown in Fig. 2.
In this diagram every chain giving essential IR divergences contains the factor
$$\begin{array}{c}\text{ }\mathrm{\Delta }_{+\rho }G^{\rho \nu }\mathrm{\Delta }_{\nu \alpha },(\mathrm{A2}.6)\text{ }\hfill \end{array}$$
where $`G^{\rho \nu }`$ is the nearest to the separated part nonseparated one which is one-particle irreducible subdiagram and $`\mathrm{\Delta }`$ are the propagators of the lines entering to the $`G^{\rho \nu }`$. The divergence takes place only when the index $`\rho `$ in (Appendix 2) takes values $`1,2`$.
Second case. We now assume that a part or all of the four-dimensional momenta lend to certain finite values. The consideration is similar. The gluon propagator now gives a pole of the second order, whereas the fermion propagator gives no poles, because of the presence of the mass. Notice, that at the investigation of IR divergences one can replace all fermion lines by $`1/M`$ which makes impossible getting $`1/M`$ to the numerator. Besides in this case one can replace $`R`$ by 1 and $`\widehat{\mu }^2`$ by $`\mu ^2`$ in the propagator of gluon which excludes $`\mathrm{\Lambda }`$ from the consideration. That is why from dimensional considerations, it follows that every nonseparated part contributes no less than four minus the number of its external lines to $`\sigma `$. This is also true for the divergent diagrams with the stipulation that for the two-point diagram, we must subtract the part independent of the momentum and proportional to $`g^{\mu \nu }`$ together with the divergent part, i.e., we must normalize the gluon mass to zero in every order. The latter does not violate the Ward identities in the limit of removing the regularization. Therefore, we now have
$$\begin{array}{c}\text{ }m4n(2sr),\text{ }\hfill \\ \text{ }\sigma =4(sn)+m2s,(\mathrm{A2}.7)\text{ }\hfill \end{array}$$
whence we obtain
$$\begin{array}{c}\text{ }\sigma r>0,(\mathrm{A2}.8)\text{ }\hfill \end{array}$$
i.e., there is no divergence.
## Appendix 3
In the framework of the perturbation theory constructed based on Lagrangian (1. Introduction)-(1. Introduction), we consider an arbitrary diagram with three gluon external lines and investigate its structure with respect to the labels of the global gauge group. This structure has the form $`A^{abc}`$. The gluon propagator is proportional to $`\delta ^{ab}`$, the three-point gluon vertex is proportional to $`f^{abc}`$, and the four-point vertex is proportional to $`f^{abe}f^{cde}`$ or $`\delta ^{ab}\delta ^{cd}`$. In addition, there can be fermion loops.
Every fermion loop gives the factor
$$\begin{array}{c}\text{ }\mathrm{ReTr}(\lambda ^a\mathrm{}\lambda ^c),(\mathrm{A3}.1)\text{ }\hfill \end{array}$$
if the number of vertices in the loop is even or
$$\begin{array}{c}\text{ }i\mathrm{ImTr}(\lambda ^a\mathrm{}\lambda ^c),(\mathrm{A3}.2)\text{ }\hfill \end{array}$$
if this number is odd (similar to the Furry theorem in quantum electrodynamics). We consider only the factors related to the gauge labels. From (1. Introduction), it follows that
$$\begin{array}{c}\text{ }\lambda _{\alpha \beta }^a\lambda _{\gamma \delta }^a=2\left(\delta _{\alpha \delta }\delta _{\beta \gamma }\frac{1}{3}\delta _{\alpha \beta }\delta _{\gamma \delta }\right).(\mathrm{A3}.3)\text{ }\hfill \end{array}$$
We can also write
$$\begin{array}{c}\text{ }\mathrm{Tr}\left(\lambda ^a\lambda ^b\lambda ^c\right)=i2f^{abc}+2d^{abc},(\mathrm{A3}.4)\text{ }\hfill \end{array}$$
whence it follows that
$$\begin{array}{c}\text{ }f^{abc}=\frac{i}{4}\left(\mathrm{Tr}\left(\lambda ^a\lambda ^b\lambda ^c\right)\mathrm{Tr}\left(\lambda ^c\lambda ^b\lambda ^a\right)\right).(\mathrm{A3}.5)\text{ }\hfill \end{array}$$
We replace all $`f^{abc}`$ in the diagram with the right-hand side of formula (Appendix 3) and then replace all resulting expressions of the form $`\lambda ^a\lambda ^a`$ according to formula (Appendix 3). As a result, only three matrices $`\lambda `$ carrying three external labels of the diagram remain. These matrices are connected with each other by their labels of the fundamental representation (because after formula (Appendix 3) is used, only the Kronecker symbols remain), i.e., $`A^{abc}`$ consists of terms of the form
$$\begin{array}{c}\text{ }\mathrm{Tr}(\lambda ^a\lambda ^b\lambda ^c).(\mathrm{A3}.6)\text{ }\hfill \end{array}$$
After using formula (Appendix 3), we have the factor $`i^N`$, where $`N`$ is the number of three-point gluon vertices. The use of formula (Appendix 3) gives no new imaginary factors. The initial expression consists of the real quantities $`\delta ^{ab}`$ and $`f^{abc}`$ and the quantities (Appendix 3) and (Appendix 3), of which the first is real and the second is imaginary; therefore, it is proportional to $`i^N^{}`$, where $`N^{}`$ is the number of fermion loops with an odd number of vertices. It is easy to show that for the diagram with three gluon external lines, the quantities $`N`$ and $`N^{}`$ have different parities. Therefore, for the powers of $`i`$ to be consistent, it is necessary that the sum of the expressions of form (Appendix 3) be imaginary, whence, using (Appendix 3), we conclude that
$$\begin{array}{c}\text{ }A^{abc}f^{abs}.(\mathrm{A3}.7)\text{ }\hfill \end{array}$$
We now similarly consider the structure of an arbitrary diagram with four gluon external lines. It has the form $`A^{abcd}`$. The quantity $`A^{abcd}`$ consists of terms of the forms $`\mathrm{Tr}(\lambda ^a\lambda ^b\lambda ^c\lambda ^d)`$ and $`\mathrm{Tr}(\lambda ^a\lambda ^b)\mathrm{Tr}(\lambda ^c\lambda ^d)`$ (the latter expression is explicitly real). It is easy to show that for the diagram with four gluon external lines, the quantities $`N`$ and $`N^{}`$ have the same parities. Therefore, the sum of the indicated expressions, which compose $`A^{abcd}`$, must be real. We can represent the quantity $`\mathrm{Tr}(\lambda ^a\lambda ^b\lambda ^c\lambda ^d)`$) as a sum of its symmetrized part (explicitly real) and expressions of the form $`if^{abe}\mathrm{Tr}(\lambda ^c\lambda ^d\lambda ^e)`$. In view of the reality condition, the latter expression must be a sum of quantities of the form $`f^{abe}f^{cde}`$. Direct calculation shows that the symmetrized part of the quantity $`\mathrm{Tr}(\lambda ^a\lambda ^b\lambda ^c\lambda ^d)`$ is proportional to the symmetrized part of the quantity $`\delta ^{ab}\delta ^{cd}`$. We can therefore conclude that $`A^{abcd}`$ consists of terms of the forms
$$\begin{array}{c}\text{ }f^{abc}f^{abc}\mathrm{and}\delta ^{ab}\delta ^{cd}.(\mathrm{A3}.8)\text{ }\hfill \end{array}$$
## Appendix 4
We shall find a form of the values which are necessary to add to Feynman diagrams in order to in the limit $`\mathrm{\Lambda }\mathrm{}`$ make them finite and coincident with its value calculated using dimensional regularization. We consider a diagram in Euclidean space constructed from propagator (3. Comparison of the light-front and Lorentz-invariant perturbation theories), from the propagator $`\mathrm{\Delta }^\psi `$, and from the vertices entering (3. Comparison of the light-front and Lorentz-invariant perturbation theories) with restrictions (1. Introduction) removed. The other vertices entering (1. Introduction) are involved only in reducing the subdiagrams of the preceding orders of the perturbation theory to the ”correct” (i.e., calculated using dimensional regularization) value including cancellation of divergences.
It is clear that the summary fermion propagator can be represented as
$$\begin{array}{c}\text{ }\mathrm{\Delta }^\psi =\frac{k_\mu \gamma ^\mu }{k^2M_0^2+i0}R^{}+\frac{M_0}{k^2M_0^2+i0}R^{\prime \prime },(\mathrm{A4}.1)\text{ }\hfill \end{array}$$
where $`R^{}`$ and $`R^{\prime \prime }`$ are cutoff factors that properly decrease and $`R^{},R^{\prime \prime }\underset{\mathrm{\Lambda }\mathrm{}}{}1`$. We represent the diagram as a sum such that one summand contains only one term of the numerator of every fermion propagator. Each of these summands can be represented as
$$\begin{array}{c}\text{ }I=𝑑kF(k,p)f_\mathrm{\Lambda }(k,p),(\mathrm{A4}.2)\text{ }\hfill \end{array}$$
where $`dk`$ represents all volume elements, $`f_\mathrm{\Lambda }(k,p)`$ is the product of all factors $`R^{}`$ and $`R^{\prime \prime }`$ entering tlie fermion propagators and of all factors $`R`$ entering boson propagators (3. Comparison of the light-front and Lorentz-invariant perturbation theories), $`k`$ denotes the integration momenta, and $`p`$ denotes the external momenta of the diagram.
It is evident that we can find a function $`\widehat{F}(k,p)`$ that is polynomial in $`p`$, is Lorentz invariant, and has no nonintegrable IR singularities such that for $`\widehat{F}(k,p)`$ and $`F(k,p)`$, a number of the first terms of their asymptotic expansions at infinity with respect to $`k`$ coincide and the difference $`F(k,p)\widehat{F}(k,p)`$ is integrable (see the refinement below). Therefore, we can write
$$\begin{array}{c}\text{ }I=𝑑k(F(k,p)\widehat{F}(k,p))f_\mathrm{\Lambda }(k,p)+𝑑k\widehat{F}(k,p)f_\mathrm{\Lambda }(k,p).(\mathrm{A4}.3)\text{ }\hfill \end{array}$$
Up to corrections of the order $`1/\mathrm{\Lambda }`$, we can neglect the factor $`f_\mathrm{\Lambda }(k,p)`$ in the first integral (because the integral converges even in the absence of this factor). We can then assume that this integral is calculated using dimensional regularization and we can split it into two integrals (assuming that they are both renormalized by dimensional regularization). As a result, we obtain the expression
$$\begin{array}{c}\text{ }I=^{dim}𝑑kF(k,p)^{dim}𝑑k\widehat{F}(k,p)+𝑑k\widehat{F}(k,p)f_\mathrm{\Lambda }(k,p).(\mathrm{A4}.4)\text{ }\hfill \end{array}$$
Using the expansion for the function $`f_\mathrm{\Lambda }(k,p)`$ in the last integral in (Appendix 4) with respect to $`p`$ in the neighborhood of the origin (this expansion is well defined for any $`k`$), we can now show that the last integral in (Appendix 4) is a polynomial in $`p`$ plus corrections of the order $`1/\mathrm{\Lambda }`$.
We must make the following refinement of these considerations. Before subtracting the overall divergence of the diagram, we must verify that the diagram has no subdivergences with respect to a part of the integration variables. Such subdivergences can be produced by subdiagrams (which is taken into account by the renormalization in the lower orders) or by the divergence on integrating over a part of the momentum components. For the QCD in the gauge $`A_{}=0`$, the latter is possible because of the improper decrease of the propagator in the direction $`k_{}`$. It is known that this divergence is present if the index of the UV divergence with respect to a part of the components is nonnegative. From the structure of the propagators, we can see that in the Euclidean space, the subtraction of the overall divergence cannot decrease only the index $`\omega _{}`$ of the UV divergence with respect to $`k_{}`$. This results in the necessity to preliminarily subtract, the subdivergence with respect to $`k_{}`$, the subdivergence being in general dependent on the projections $`p^\nu n_\mu `$ of the external momenta in an arbitrary nonpolynomial way. Analysis of the diagrams for the Green’s functions shows that we have $`\omega _{}0`$ for only the one-loop two-point diagrams and that $`\omega _{}=0`$. These diagrams have only one external momentum. Therefore, they cannot have a nonpolynomial dependence on $`p^\nu n_\nu `$, because of the invariance with respect to multiplying the vector $`n_\nu `$ by a complex number. One must take into account that, for the diagrams of the Green’s functions, the vector $`n_\nu `$ with the nonconvolute index $`\nu `$ cannot stand in the numerator, because this gives zero on convolution with the propagators. If we consider the one-particle irreducible vertex parts, whose diagrams do not satisfy the last condition, the divergent parts can be nonpolynomial , and the number of diagrams that are divergent with respect to $`k_{}`$ can be much larger.
It seems that if the integral of the form, described at the beginning of this Appendix, converges then the result of it’ s calculation cannot depend on the vector $`n_\nu ^{}`$ (complex conjugated to $`n_\nu `$ up to a factor). But this would be true only if this integral be an analitical function of the complex vector $`n_\nu `$. However the derivative of our integral with respect to complex vector $`n_\nu `$ can be nonconvergent due to the rising of infrared singularity (of the power of the pole in $`(nk)`$). Therefore the result of the calculation of the diagram, and, hence, it’s divergent part can depend on $`n_\nu `$ and on $`n_\nu ^{}`$. It is seen from the equation (3. Comparison of the light-front and Lorentz-invariant perturbation theories) for the propagator that the integrands we considered are invariant with respect to a multiplication of the vector $`n_\nu `$ by a complex number (and of the $`n_\nu ^{}`$ by complex conjugate number). This allows to conclude that up to corrections of the order $`1/\mathrm{\Lambda }`$ the difference between the diagram described in the beginning of this appendix (more precisely, the finite sum of such diagrams of the given order) and its value calculated using dimensional regularization (i.e., the similar sum of the first integrals in (Appendix 4)) reduces to a polynomial in the external momenta with coefficients containing factors $`N_{\alpha \beta }`$ (see the definition of this value in formula (1. Introduction)). |
warning/0002/cond-mat0002155.html | ar5iv | text | # Solution of the Schrödinger Equation for Quantum Dot Lattices with Coulomb Interaction between the Dots
## I Introduction
Quantum dots have been in the focus of intensive research already for at least a decade which lead to a countless number of publications <sup>*</sup><sup>*</sup>*Therefore we will refer here only to papers which are directly connected to the scope of this work (for a recent book see Ref.). Although almost all experiments are performed at dot lattices, in the vast majority of theoretical investigations the interaction between dots is neglected. This is for the following reasons: i) Because the confinement frequency $`\omega _0`$ is a parameter, which is mainly extracted from optical properties, it is difficult to tell the influence of dot interaction apart from the intrinsic single– dot value. (Possibilities to overcome this problem are discussed in the present work.) ii) The theory of Raman spectra, which can in principle monitor the dispersion (wave number dependence) of excitation energies as a direct consequence of interdot– interaction, is not yet advanced enough to extract the dispersion. iii) The lattice constant of dot arrays produced with current technologies is so large ($`>2000`$Å) that large electron numbers $`N`$ per dot are necessary to obtain a seizable amount of shift. For these N, however, reliable first principle calculations are not possible. With the advent of self– assembled dot arrays the last item might change.
The scope of this paper is to investigate conditions, which lead to qualitative and observable effects of interdot– interaction on excitation spectra and the phase transition found in Ref.. Unlike in Ref., a magnetic field B is explicitly taken into account and a microscopic theory is applied. Our approach is purely microscopic, i.e. we solve the Schrödinger equation of a model system exactly. Our model comprises the following approximations: i) The dot confinement is strictly parabolic in radial direction, but with anisotropic confinement frequencies $`\omega _i(i=1,2)`$ and independent of $`N`$ and B. ii) Overlap of wave functions between different dots is neglected (no hopping). iii) The Coulomb interaction of the electrons in different dots is treated in dipole approximation (second order in dot diameter over lattice constant). Our model is similar to that in Ref., but allows more complicated dots and lattice structures. Besides we calculate also the intra dot excitations (apart from the collective center– of– mass excitations) for $`N=2`$ explicitly and discuss the instability in this microscopic model. Our results on the lateral dot dimer are compared with a former paper , which uses a high magnetic field approach, in Sect. III.
The plan of this paper is as follows. For further reference, we briefly summarize in Sect. I some relevant results for one single dot, or for dot lattices, where the distance between the dots is very large. This is important, because all exact solutions in the center– of– mass subsystem are traced back (by special transformations) to the solution of this one– electron Hamiltonian. This is analogous to ordinary molecular and lattice dynamics. After this, we consider a dot dimer, which mimics a lattice, where the dots are pairwise close to each other. This model can give an idea of the effects expected in dot lattices with a basis. Next we consider a rectangular, but primitive lattice in order to obtain the dispersion in the spectra. Finally, the intra-dot excitations of the Hamiltonian in the relative coordinates are calculated numerically. The paper ends with a summary. In the Appendix we give a short and elementary proof for the fact that the Generalized Kohn Theorem holds even for arbitrary arrays of identical non-circular quantum dots with Coulomb interaction (between the dots) in an homogeneous magnetic field.
## II Single Dot
The Hamiltonian considered here reads (in atomic units $`\mathrm{}=m=e=1`$)
$$H=\underset{i=1}{\overset{N}{}}\left\{\frac{1}{2m^{}}\left[𝐩_i+\frac{1}{c}𝐀(𝐫_i)\right]^2+\frac{1}{2}𝐫_i𝐂𝐫_i\right\}+\frac{1}{2}\underset{ik}{}\frac{\beta }{|𝐫_i𝐫_k|}$$
(1)
where $`m^{}`$ is the effective mass (in units of the bare electron mass $`m`$), $`\beta `$ the inverse dielectric constant of the background, and $`𝐂`$ a symmetric tensor. In case of a single dot, $`𝐂`$ is given by the confinement potential and we define $`𝐂=𝛀`$. It is always possible to find a coordinate system where $`\mathrm{\Omega }_{12}=\mathrm{\Omega }_{21}=0`$ and $`\mathrm{\Omega }_{ii}=\omega _i^2=m^{}\omega _i^2`$. We use the symmetric gauge $`𝐀=\frac{1}{2}𝐁\times 𝐫`$ throughout. The Zeeman term in $`H`$ is disregarded at the moment.
For $`N=1`$, the Hamiltonian
$$H=\frac{1}{2m^{}}\left[𝐩+\frac{1}{c}𝐀(𝐫)\right]^2+\frac{1}{2}𝐫𝐂𝐫$$
(2)
can be diagonalized exactly. Later on we will see that also the case of interacting dots can be traced back to the solution of type (2). (Therefore, we kept the off diagonal elements of $`𝐂`$ in the results given below because the dynamical matrix, which also contributes to $`𝐂`$, is generally non–diagonal and we want to use the same coordinate system for all $`𝐪`$ values.) After transforming the operators $`𝐫_i`$ and $`𝐩_i`$ to creation– annihilation operators (see e.g. Ref.) and using the procedure described by Tsallis , we obtain for the eigenvalues
$$E(n_+,n_{})=(n_++\frac{1}{2})\omega _++(n_{}+\frac{1}{2})\omega _{};n_\pm =0,1,2,\mathrm{}$$
(3)
where
$`\omega _\pm `$ $`=`$ $`\sqrt{{\displaystyle \frac{\omega _c^2}{2}}+\stackrel{~}{\omega }_0^2\pm \sqrt{{\displaystyle \frac{\omega _c^4}{4}}+\omega _c^2\stackrel{~}{\omega }_0^2+{\displaystyle \frac{\mathrm{\Delta }^2}{4}}+C_{12}^2}}`$ (4)
$`=`$ $`\sqrt{\left[{\displaystyle \frac{1}{2}}\sqrt{\omega _c^2+4\stackrel{~}{\omega }_0^2+{\displaystyle \frac{(\mathrm{\Delta }^2+4C_{12}^2)}{\omega _c^2}}}\pm {\displaystyle \frac{\omega _c^{}}{2}}\right]^2{\displaystyle \frac{(\mathrm{\Delta }^2+4C_{12}^2)}{4\omega _c^2}}}`$ (5)
$$\stackrel{~}{\omega }_0^2=\frac{1}{2}(C_{11}+C_{22});\mathrm{\Delta }=C_{11}C_{22}$$
(6)
and $`\omega _c^{}=\frac{B}{m^{}c}`$ is the cyclotron frequency with the effective mass. (The results for the special case $`C_{12}=0`$ can also be found in Ref. .) The optical selection rules are the same as in the circular case, i.e., there are two possible types of excitations
$$(\mathrm{\Delta }n_+=\pm 1\text{and}\mathrm{\Delta }n_{}=0)\text{or}(\mathrm{\Delta }n_{}=\pm 1\text{and}\mathrm{\Delta }n_+=0)$$
(7)
leading to the excitation energies $`\mathrm{\Delta }E=\omega _+`$ and $`\omega _{}`$. It is easily seen that the form (5) reduces to the familiar formula in the circular case, where $`\mathrm{\Delta }=0`$ and $`C_{12}=0`$. By inspection of (4) we find that a soft mode $`\omega _{}(B)=0`$ can only occur if $`C_{11}C_{22}=C_{12}^2`$. For a diagonal $`𝐂`$ this means that $`min(C_{11},C_{22})=0`$. The last condition is of importance for interacting dots considered in the next Sections.
In the limiting case $`B=0`$ we obtain from (4)
$$\omega _\pm (B=0)=\sqrt{\frac{(C_{11}+C_{22})}{2}\pm \sqrt{\frac{(C_{11}C_{22})^2}{4}+C_{12}^2}}$$
(8)
We see that degeneracy $`\omega _+(B=0)=\omega _{}(B=0)`$ can only happen if $`C_{12}=0`$ and $`C_{11}=C_{22}`$. For a diagonal confinement tensor with $`C_{12}=0`$ we obtain $`\omega _+(B=0)=\text{max}(\omega _1,\omega _2)`$ and $`\omega _{}(B=0)=\text{min}(\omega _1,\omega _2)`$. As to be expected, we observe a gap between the two excitation curves $`\omega _+(B)`$ and $`\omega _{}(B)`$ at $`B=0`$, if the two confinement frequencies do not agree.
Alternatively we can introduce the quantum numbers
$$k=\frac{(n_++n_{})|n_++n_{}|}{2};m_z=n_+n_{}$$
(9)
where $`k`$ is the node number and $`m_z`$ turns in the circular limit into the angular momentum quantum number.
For arbitrary $`N`$, the center of mass (c.m.) $`𝐑=\frac{1}{N}_i𝐫_i`$ can be separated $`H=H_{c.m.}+H_{rel.}`$ with
$$H_{c.m.}=\frac{1}{N}\left\{\frac{1}{2m^{}}\left[𝐏+\frac{N}{c}𝐀(𝐑)\right]^2+\frac{N^2}{2}𝐑𝐂𝐑\right\}$$
(10)
where $`𝐏=i_𝐑`$ (see Appendix). As well known, $`H_{c.m.}`$ does not contain the electron– electron– interaction. $`H_{c.m.}`$ can be obtained from the one– electron Hamiltonian (2) by the substitution: $`BNB`$$`𝐂N^2𝐂`$ and $`H\frac{1}{N}H`$. If we make the same substitution in the eigenvalues (3), we obtain
$$E_{c.m.}(n_+,n_{})=E_{N=1}(n_+,n_{})$$
i.e., the eigenvalues of the c.m. Hamiltonian are independent of $`N`$. In other words, in $`H`$ there are excitations, in which the pair correlation function is not changed, or classically speaking, where the charge distribution oscillates rigidly. Because FIR radiation (in the limit $`\lambda \mathrm{}`$) can excite only the c.m. subspace, all we see in FIR spectra is the c.m. modes.
## III Dot Dimer
We consider two identical elliptical dots centered at $`𝒂_1=(a/2,0)`$ and $`𝒂_2=(+a/2,0)`$. We expand the Coulomb interaction between electrons in different dots in a multi-pole series and restrict ourselves to the dipole approximation. By introduction of c.m. and relative coordinates within each dot, the c.m. coordinates and the relative coordinates decouple The tilde indicates that in this preliminary Hamiltonian a common gauge center for both dots is used.
$$\stackrel{~}{H}=\stackrel{~}{H}_{c.m.}(𝐑_1,𝐑_2)+\underset{\alpha }{\overset{1,2}{}}\stackrel{~}{H}_\alpha \left(\{𝐫\}_\alpha ^{(N1)}\right)$$
(11)
$`\{𝐫\}_\alpha ^{(N1)}`$ symbolizes $`(N1)`$ relative coordinates in the $`\alpha ^{th}`$ dot. This means, we have 3 decoupled Hamiltonians: the c.m. Hamiltonian and two Hamiltonians in the relative coordinates of either dot. This leads to two types of excitations:
i) Collective excitations from $`\stackrel{~}{H}_{c.m.}`$ which involve the c.m. coordinates of both dots simultaneously. Because of the harmonic form (in the dipole approximation), there are exactly two modes per dot, thus a total of four. Each excitation can be classically visualized as vibrations of rigidly moving charge distributions of both dots.
ii) Intra dot excitations which are doubly degenerate for two identical dots. Because $`\stackrel{~}{H}_\alpha (\{𝐫\}_\alpha ^{(N1)})`$ is not harmonic (it includes the exact Coulomb interaction between the electrons within each dot, which is not harmonic), this spectrum is very complex. It is the excitation spectrum of a single dot in a modified confinement potential where the c.m. coordinate is fixed. The extra term in the modified confinement potential comes from the dipole contribution of the interdot Coulomb interaction.
In this Section we consider only the c.m. Hamiltonian and focus our attention to the the effects of ellipticity in the dot confinement potential. The relative Hamiltonian for $`N=2`$ is explicitly given in the last Section and solved for circular dots. For the elliptical confinement potential considered in this Section, the relative Hamiltonian cannot be solved easily, even if we restrict ourselves to $`N=2`$, because the elliptic confinement potential breaks the circular symmetry of the rest of the relative Hamiltonian.
### A Center of Mass Hamiltonian of the Dimer
The c.m. Hamiltonian in the dipole approximation reads
$$\stackrel{~}{H}_{c.m.}=\frac{1}{N}\left\{\underset{\alpha }{\overset{1,2}{}}\frac{1}{2m^{}}\left[𝐏_\alpha +\frac{N}{c}𝐀(𝐔_\alpha +𝒂_\alpha )\right]^2+\frac{N^2}{2}\underset{\alpha ,\alpha ^{}}{}𝐔_\alpha 𝐂_{\alpha ,\alpha ^{}}𝐔_\alpha ^{}\right\}$$
(12)
where the small elongation $`𝐔_\alpha `$ is defined by $`𝐑_\alpha =𝒂_\alpha +𝐔_\alpha `$ and $`𝐏=i_𝐑=i_𝐔`$. The tensor $`𝐂`$ is
$`C_{\alpha ,\alpha }`$ $`=`$ $`𝛀+\beta N{\displaystyle \underset{\alpha ^{}(\alpha )}{}}𝐓(𝒂_{\alpha ,\alpha ^{}})`$ (13)
$`C_{\alpha ,\alpha ^{}}`$ $`=`$ $`\beta N𝐓(𝒂_{\alpha ,\alpha ^{}})\text{for}\alpha \alpha ^{}`$ (14)
where $`𝒂_{\alpha ,\alpha ^{}}=𝒂_\alpha 𝒂_\alpha ^{}`$ , and the dipole tensor is
$$𝐓(𝒂)=\frac{1}{a^5}\left[3𝒂𝒂a^2𝐈\right]$$
(15)
containing a dyad product ($``$) and the unit tensor $`𝐈`$. As in the c.m. system of a single dot, the explicit $`N`$– dependence in (12) cancels in the eigenvalues. What is left is only the $`N`$– dependence in the dipole contribution of the dot interaction appearing in (13) and (14). This means, that the c.m. spectrum of interacting dots is no longer independent of $`N`$.
The term $`𝒂_\alpha `$ in the argument of the vector potential in (12) causes trouble in finding the eigenvalues. This shift is a consequence of the fact that we have to adopt a common gauge center for both dots (we chose the middle between both dots). This problem can be solved by applying the following unitary transformation
$$H_{c.m.}=Q^1\stackrel{~}{H}_{c.m.}Q;Q=\underset{\alpha }{\overset{1,2}{}}e^{i\frac{N}{2c}(𝐁\times 𝒂_\alpha )𝐔_\alpha }$$
(16)
In other words, $`H_{c.m.}`$ agrees with $`\stackrel{~}{H}_{c.m.}`$ except for the missing shift in the argument of the vector potential.
The 4 modes inherent in $`H_{c.m.}`$ are not yet explicitly known, because the 4 degrees of freedom are coupled. Decoupling into two oscillator problems of type (2) can be achieved by the following transformation:
$$𝐔^{(+)}=\frac{1}{2}(𝐔_2+𝐔_1);𝐔^{()}=𝐔_2𝐔_1$$
(17)
This results in
$$H_{c.m.}=\frac{1}{2}H^{(+)}+\mathrm{\hspace{0.33em}2}H^{()}$$
(18)
where
$`H^{(+)}`$ $`=`$ $`{\displaystyle \frac{1}{N}}\left\{{\displaystyle \frac{1}{2m^{}}}\left[𝐏^{(+)}+{\displaystyle \frac{2N}{c}}𝐀\left(𝐔^{(+)}\right)\right]^2+{\displaystyle \frac{N^2}{2}}𝐔^{(+)}\left(4𝛀\right)𝐔^{(+)}\right\}`$ (19)
$`H^{()}`$ $`=`$ $`{\displaystyle \frac{1}{N}}\{{\displaystyle \frac{1}{2m^{}}}[𝐏^{()}+{\displaystyle \frac{N}{2c}}𝐀\left(𝐔^{()}\right)]^2+{\displaystyle \frac{N^2}{2}}𝐔^{()}({\displaystyle \frac{1}{4}}𝛀+{\displaystyle \frac{N}{2}}\beta 𝐓(𝒂))𝐔^{()})\}`$ (20)
and $`𝒂`$ is a vector pointing from one dot center to the other. Then $`𝐓(𝒂)`$ has the following components
$$T_{11}=\frac{2}{a^3};T_{22}=\frac{1}{a^3};T_{12}=T_{21}=0$$
(21)
Now we assume that the principle axes of the confinement potentials are in x-y-direction. This means
$$\mathrm{\Omega }_{11}=\omega _1^2;\mathrm{\Omega }_{22}=\omega _2^2;\mathrm{\Omega }_{12}=\mathrm{\Omega }_{21}=0$$
(22)
The eigenvalues of $`H^{(+)}`$ can be obtained from (3) and (4) with
$$\stackrel{~}{\omega }_0^2=\frac{1}{2}(\omega _1^2+\omega _2^2);\mathrm{\Delta }=(\omega _1^2\omega _2^2)$$
(23)
and for $`H^{()}`$ with
$$\stackrel{~}{\omega }_0^2=\frac{1}{2}(\omega _1^2+\omega _2^2)+\frac{1}{2}p;\mathrm{\Delta }=(\omega _1^2\omega _2^2)+3p$$
(24)
where the interaction parameter is defined by
$$p=\frac{2N\beta }{a^3}$$
(25)
Observe that the dependence on $`N`$ cancels, except that included in $`p`$ (see discussion following (10)).
It is important that the dot interaction influences the result only through a single parameter. This conclusion agrees with the semi– phenomenological theory in Ref..
In all our figures we express frequencies in units of the average confinement frequency $`\omega _0=\frac{1}{2}(\omega _1+\omega _2)`$, and $`\mathrm{\Delta }`$ and $`p`$ in units $`\omega _0^2`$. Then, all systems can be characterized by the two parameters: $`\omega _1/\omega _2`$ and $`p`$. In other words, all systems having the $`\omega _1/\omega _2`$ ratio indicated in the figures are represented by the family of curves with the $`p`$ values shown. The only exception we made is the cyclotron frequency $`\omega _c^{}`$. $`\omega _c^{}/\omega _0`$ would be a good parameter in this sense, but we chose to use the magnetic field in $`Tesla`$ instead for better physical intuition. The conversion between both scales is given by $`\omega _c^{}[a.u.^{}]=\frac{0.913410^2}{m^{}}B[Tesla]`$ or, $`\omega _c^{}[\omega _0]=\frac{0.913410^2}{m^{}\omega _0[a.u.^{}]}B[Tesla]`$ In this paper we used $`\omega _0=0.2a.u.^{}=2.53meV`$ and $`m^{}`$ of GaAs. (We want to stress that this choice effects only the magnetic field scale and not the qualitative features of the figures.) For easy comparison with experimental parameters we add the definition of effective atomic units ($`a.u.^{}`$) in GaAs ($`m^{}=0.067,\beta =1/12`$) for the energy: $`1a.u.^{}=4.6510^4doubleRydberg=12.64meV`$ , and for the length: $`1a.u.^{}=1.79110^2Bohrradii=0.947710^2\AA `$.
Because $`𝐔^{(+)}`$ agrees with the total c.m. $`𝐑=\frac{1}{2}(𝐑_1+𝐑_2)`$ of the system, $`H^{(+)}`$ is the total c.m. Hamiltonian. For $`B=0`$, the eigenmodes can be visualized by classical oscillations. The two eigenmodes of $`H^{(+)}`$ are (rigid) in– phase oscillations of the dots in x and y direction, respectively. Because of the Kohn theorem (see Appendix), the independence of $`H^{(+)}`$ on the Coulomb interaction does not only hold in the dipole approximation, but it is rigorous. This shows also that the dipole approximation is consistent with the Kohn theorem, which is not guaranteed for single particle approaches. Because FIR radiation excites (in the dipole approximation) only the c.m. modes, it is only the $`p`$–independent eigenmodes of $`H^{(+)}`$ which are seen in FIR absorption experiments.
This statement is in contradiction to Ref.. They performed numerical diagonalizations for a lateral pair of circular dots confining the set of basis functions to the lowest Landau level and considering parallel spin configurations only. This is justified in the limit of high magnetic fields. They found a splitting of the two dipole allowed modes at $`B=0`$ due to dot interaction and some anti-crossing structures in the upper mode, whereas the lower mode is always close to the single particle mode. This fact is already a strong indication that the missing higher Landau levels cause both spurious effects. (Observe that the lifting of the degeneracy at $`B=0`$ in the dipole allowed excitations in Fig.1 is due to the ellipticity of the intrinsic confinement and not due to dot interaction.)
The eigenvalues of $`H^{()}`$ do depend on $`p`$ because the dots oscillate (rigidly) in its two eigenmodes in opposite phase, one mode in $`x`$ and one mode in $`y`$ direction. This leads to a change in the Coulomb energy. The two eigenmodes of $`H^{()}`$ can also be described as a breathing mode (in x direction) and a shear mode (in y direction).
### B Special Features of the Excitation Spectrum
In Fig.1a and 1b, the four excitation frequencies of the dimer are shown with $`p`$ as a parameter.
For $`p=0`$, the two modes $`\omega _\pm ^{()}`$ agree with the two modes $`\omega _\pm ^{(+)}`$. In all symbols, the superscript sign refers to the system $`H^{(+)}`$ and $`H^{()}`$ (c.m. or relative coordinate), and the subscript sign discriminates the two modes of the same system. The two modes $`\omega _\pm ^{(+)}`$ are independent of $`p`$. There are two qualitatively different cases. (Consider that $`\omega _1`$ is the oscillator frequency parallel to the line, which connects the two dot centers, and $`\omega _2`$ is the oscillator frequency perpendicular to this line.) If $`\omega _1\omega _2`$ (Fig.1a), the gap between $`\omega _+^{()}`$ and $`\omega _{}^{()}`$ at $`B=0`$ increases steadily with increasing $`p`$ until, for a critical $`p_{cr}=\omega _2^2`$ (in our numerical case: $`p_{cr}[\omega _0^2]=16/25=0.64`$) the lower mode $`\omega _{}^{()}`$ becomes soft. This transition is independent of $`B`$. For $`\omega _1\omega _2`$ the gap between $`\omega _+^{()}`$ and $`\omega _{}^{()}`$ at $`B=0`$ first decreases with increasing $`p`$ until it vanishes for $`p=\frac{1}{3}(\omega _2^2\omega _1^2)`$ (in our numerical case: $`p[\omega _0^2]=4/15=0.27`$). Afterwards, it increases until the lattice becomes soft at $`p_{cr}=\omega _2^2`$ (in our numerical case: $`p_{cr}[\omega _0^2]=36/25=1.44`$). The dependence of the two excitation energies $`\omega _+^{()}`$ and $`\omega _{}^{()}`$ on $`p`$ for $`B=0`$ in the second case is shown in Fig.2. Comparison of Fig.s 2a and 2b demonstrates that the dot architecture in Fig.2a is much more sensitive to interdot interaction than that in Fig.2b. Thus, if we want to observe or use the instability, this event happens in case 2a for for smaller $`p`$ (or larger lattice constants) than in case 2b. Additionally, the assumption of non–overlapping dot wave functions (for a given lattice constant) is better fulfilled in case 2a than in case 2b.
For GaAs as a typical substance, (25) can be rewritten in more convenient units as
$$p[\omega _0^2]=\frac{2.2610^7N}{(a[\AA ])^3\left(\omega _0[meV]\right)^2}$$
(26)
Obviously, we need large dots (large $`N`$, small $`\omega _0`$– which means large polarizability), and a small dot distance $`a`$ for a seizable interaction effect. On the other hand, the dot radius for $`N=1`$ is of order of the effective magnetic length $`l_0=\left((2\omega _0)^2+(\omega _c^{})^2\right)^{1/4}`$, which reads for GaAs
$$l_0[\AA ]=\frac{238}{\left(\left(\omega _0[meV]\right)^2+0.739\left(B[Tesla]\right)^2\right)^{\frac{1}{4}}}$$
(27)
and we need small dots and high magnetic fields for small overlap. Consequently, a magnetic field helps avoiding overlap of the dots, although e.g. the critical $`p`$ for soft modes is independent of $`B`$. The question is, if there exists a window between these two (partly) conflicting demands. For an order– of– maigntude estimate, let us consider GaAs with $`\omega _0`$ as chosen above and the worst case $`N=1`$. Then (26) with a typical $`p[\omega _0^2]=0.1`$ (which seems to be the minimum for any observable effect) provides a dot distance of $`a[\AA ]=327`$ and (27) gives for $`B=0`$ a radius of $`l_0[\AA ]=150`$ and for $`B[Tesla]=10`$ a radius of $`l_0[\AA ]=80`$. Consequently, the constraint $`l_0<a/2`$ for our model can be fulfilled. For obtaining larger interaction effects the parameters have to be optimized.
The next question is what happens in mode softening physically? Firstly, it is the antisymmetric shear mode $`\omega _{}^{()}`$ which has the lowest frequency and which becomes soft. If the interaction parameter is strong enough ($`p>p_{cr}`$), the decrease in interdot– Coulomb energy with increasing elongation of the dots becomes larger than the increase of confinement potential energy. Because in the harmonic model both energies depend quadratically on elongation, the dimer would be ionized, i.e. stripped of the electrons. Clearly, in this case we have to go beyond the dipole approximation for the interdot interaction and beyond the harmonic approximation for the confinement potential.
In order to obtain a hand-waving picture of what happens, the confinement potential of the system for shear mode oscillations is supplemented by a $`4^{th}`$ order term in the following way: $`V_{conf.}=2N\left[\frac{1}{2}\omega _2^2U^2AU^4\right]`$ with $`(A>0)`$, and the Coulomb interaction in $`4^{th}`$ order reads: $`V_{int}=pNU^2+(3pN/a^2)U^4`$ where $`p=2N\beta /a^3`$ as above. Then, the stability condition reads
$$\frac{V_{tot}}{N}=(\omega _2^2p)U^2+(\frac{3p}{a^2}2A)U^40$$
(28)
The condition for the existence of a bound state is that the $`U^4`$\- term is positive: $`3p/a^2>2A`$. For a positive $`U^2`$\- term ( $`p<\omega _2^2`$ ), the equilibrium position is $`U_0=0`$. If the $`U^2`$\- term becomes negative ( $`p>\omega _2^2`$ ), the system finds a new equilibrium at a finite elongation
$$U_0=\pm \sqrt{\frac{(p\omega _2^2)}{2(\frac{3p}{a^2}2A)}}$$
(29)
This new ground state is doubly degenerate: $`𝐔_1=(a,+U_0),𝐔_2=(+a,U_0)`$ and $`𝐔_1=(a,U_0),𝐔_2=(+a,+U_0)`$ have the same energy. In short, at $`p_{cr}=\omega _2^2`$ there is an electronic phase transition to a polarised state, where the equilibrium position of the c.m. is no more in the middle of the dots. At the end we want to stress that all these stability considerations are only valid if the confinement potential is not changed under elongation of the c.m. of the dots. Secondly, it is not rigorous to include the fourth order terms after separation of c.m. and relative coordinates, because in fourth order these two coordinates do not decouple exactly.
## IV Dot Lattice
We consider a periodic lattice of equal quantum dots at lattice sites $`𝐑_{n,\alpha }^{(0)}=𝐑_n^{(0)}+𝒂_\alpha `$. The vectors $`𝐑_n^{(0)}`$ form a Bravais lattice and $`𝒂_\alpha `$ runs over all sites within an unit cell. In developing a theory for these lattices we have to repeat all steps in Sect. II from (11) to (16) just by supplementing the index $`\alpha `$ by the index $`n`$ for the unit cell.
### A Center of Mass Hamiltonian of the Dot Lattice
The c.m. Hamiltonian in the dipole approximation then reads
$`H_{c.m.}`$ $`=`$ $`{\displaystyle \frac{1}{N}}\{{\displaystyle \underset{n,\alpha }{}}{\displaystyle \frac{1}{2m^{}}}[𝐏_{n,\alpha }+{\displaystyle \frac{N}{c}}𝐀\left(𝐔_{n,\alpha }\right)]^2`$ (31)
$`+{\displaystyle \frac{N^2}{2}}{\displaystyle \underset{\genfrac{}{}{0pt}{}{n,\alpha }{n^{},\alpha ^{}}}{}}𝐔_{n,\alpha }𝐂_{n,\alpha ;n^{},\alpha ^{}}𝐔_{n^{},\alpha ^{}}\}`$
where $`𝐔_{n,\alpha }=𝐑_{n,\alpha }𝐑_{n,\alpha }^{(0)}`$ is the elongation of the c.m. at lattice site $`(n,\alpha )`$ and the force constant tensor $`𝐂`$ is defined in analogy to (13) and (14).
The Hamiltonian (31) is a phonon Hamiltonian in an additional homogeneous magnetic field. The first stage of decoupling can be achieved by the usual phonon transformation
$`𝐔_{n,\alpha }`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{N_c}}}{\displaystyle \underset{𝐪}{\overset{BZ}{}}}e^{i𝐪R_n^{(0)}}𝐔_{𝐪,\alpha }`$ (32)
$`𝐏_{n,\alpha }`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{N_c}}}{\displaystyle \underset{𝐪}{\overset{BZ}{}}}e^{+i𝐪R_n^{(0)}}𝐏_{𝐪,\alpha }`$ (33)
where $`N_c`$ is the number of unit cells and the transformed coordinates have the following properties $`𝐔_{𝐪,\alpha }=𝐔_{𝐪,\alpha }^{}=𝐔_{𝐪,\alpha }^{}`$ and $`𝐏_{𝐪,\alpha }=𝐏_{𝐪,\alpha }^{}`$. The Hamiltonian in the new coordinates is a sum of $`N_c`$ subsystems of dimension $`2\times `$ number of dots per unit cell: $`H_{c.m.}=_𝐪H_𝐪`$, where
$`H_𝐪`$ $`=`$ $`{\displaystyle \frac{1}{N}}\{{\displaystyle \underset{\alpha }{}}{\displaystyle \frac{1}{2m^{}}}[𝐏_{𝐪,\alpha }+{\displaystyle \frac{N}{c}}𝐀(𝐔_{𝐪,\alpha }^{})]^{}[𝐏_{𝐪,\alpha }+{\displaystyle \frac{N}{c}}𝐀(𝐔_{𝐪,\alpha }^{})]`$ (35)
$`+{\displaystyle \frac{N^2}{2}}{\displaystyle \underset{\alpha ,\alpha ^{}}{}}𝐔_{𝐪,\alpha }^{}𝐂_{𝐪;\alpha ,\alpha ^{}}𝐔_{𝐪,\alpha ^{}}\}`$
The dynamical matrix is defined by
$$𝐂_{𝐪;\alpha ,\alpha ^{}}=\underset{n}{}e^{i𝐪𝐑_n^{(0)}}𝐂_{\alpha ,\alpha ^{}}\left(𝐑_n^{(0)}\right);𝐂_{\alpha ,\alpha ^{}}\left(𝐑_n^{(0)}\right)=𝐂_{n,\alpha ;\mathrm{\hspace{0.17em}0},\alpha ^{}}$$
(36)
and it is hermitean $`𝐂_{𝐪;\alpha ^{},\alpha }=𝐂_{𝐪;\alpha ,\alpha ^{}}^{}=𝐂_{𝐪;\alpha ,\alpha ^{}}`$.
Next we want to recover the limiting case considered in Sect. III. If the dots in a given unit cell are far away from those in neighboring cells, then in (36) only the term with $`𝐑_n^{(0)}=0`$ contributes, $`𝐂`$ does not depend on $`𝐪`$, consequently the index $`𝐪`$ is redundant, and (35) agrees with (12).
Our preliminary result (35) is not yet diagonal in $`\alpha ,\alpha ^{}`$. In some special cases (see e.g. two identical dots per unit cell considered in Sect. III) this can be achieved by an unitary transformation
$$𝐔_{𝐪,\alpha }=\underset{\alpha ^{}}{}Q_{𝐪;\alpha ,\alpha ^{}}\stackrel{~}{𝐔}_{𝐪,\alpha ^{}};Q_{𝐪;\alpha ^{},\alpha }^{}=Q_{𝐪;\alpha ,\alpha ^{}}^1$$
(37)
under which the one– particle term in (35) is invariant and the transformed interaction term
$$\frac{1}{2}\underset{\alpha ,\alpha ^{}}{}\stackrel{~}{𝐔}_{𝐪,\alpha }^{}\stackrel{~}{𝐂}_{𝐪;\alpha ,\alpha ^{}}\stackrel{~}{𝐔}_{𝐪,\alpha ^{}}\text{with}\stackrel{~}{𝐂}_{𝐪;\alpha ,\alpha ^{}}=\underset{\alpha _1,\alpha _2}{}Q_{\alpha ,\alpha _1}^1𝐂_{𝐪;\alpha _1,\alpha _2}Q_{\alpha _2,\alpha ^{}}$$
(38)
can be made diagonal $`\stackrel{~}{𝐂}_{𝐪;\alpha ,\alpha ^{}}=\stackrel{~}{𝐂}_{𝐪;\alpha }\delta _{\alpha ,\alpha ^{}}`$ by a proper choice of $`Q_{\alpha ,\alpha ^{}}`$. Now, (35) reads $`H_𝐪=_\alpha H_{𝐪,\alpha }`$, where
$`H_{𝐪,\alpha }`$ $`=`$ $`{\displaystyle \frac{1}{N}}\{{\displaystyle \frac{1}{2m^{}}}[\stackrel{~}{𝐏}_{𝐪,\alpha }+{\displaystyle \frac{N}{c}}𝐀(\stackrel{~}{𝐔}_{𝐪,\alpha }^{})]^{}[\stackrel{~}{𝐏}_{𝐪,\alpha }+{\displaystyle \frac{N}{c}}𝐀(\stackrel{~}{𝐔}_{𝐪,\alpha }^{})]`$ (40)
$`+{\displaystyle \frac{N^2}{2}}\stackrel{~}{𝐔}_{𝐪,\alpha }^{}\stackrel{~}{𝐂}_{𝐪;\alpha }\stackrel{~}{𝐔}_{𝐪,\alpha }\}`$
The eigenvalues of (40) can be obtained from those of (2) because corresponding quantities have the same commutation rules. Such an unitary transformation does not exist, e.g., for two different dots per cell. Then (35) has to be solved directly using the method described in Ref. .
### B Dynamical Matrix for Bravais lattices
From now on we consider Bravais lattices what means that we can forget the indices $`\alpha `$ in the first part of this Section. Then the dynamical matrix
$$𝐂_𝐪=𝛀+\beta N\underset{𝐑_n^{(0)}0}{}\left(1e^{i𝐪𝐑_n^{(0)}}\right)𝐓\left(𝐑_n^{(0)}\right)$$
(41)
is real and symmetric, but generally not diagonal, even if $`𝛀`$ is diagonal. A very important conclusion is apparent in (41). In the limit $`𝐪0`$, the inter– dot interaction (represented by $`\beta `$) has no influence on $`𝐂_𝐪`$ and therefore on the spectrum. This means, that the excitation spectrum observed by FIR spectroscopy is not influenced by inter– dot interaction and agrees with the one– electron result (as in the single dot). This statement is rigorous for parabolic confinement (see Appendix). It can also be understood intuitively, because a $`q=0`$– excitation is connected with homogeneous in– phase elongations of the dots which do not change the distance between the electrons. We want to mention that this conclusion seems to be in contradiction with the experimental work in Ref.. They found a splitting of the upper and lower excitation branch at $`B=0`$ and $`q=0`$ for circular dots in a rectangular lattice, which they interpreted within a phenomenological model of interacting dipoles as a consequence of lattice interaction. However, they use mesoscopic dots with a diameter of $`370000\AA `$ and lattice periods of $`400000`$ and $`800000\AA `$. These dots are clearly beyond our microscopic quantum mechanical model, which rests on a parabolic confinement.
For the rectangular lattices considered in our numerical examples we define $`𝐑^{(0)}=N_1a_1𝐞_1+N_2a_2𝐞_2`$ and $`𝐪=q_1\frac{2\pi }{a_1}𝐞_1+q_2\frac{2\pi }{a_2}𝐞_2`$ with the lattice constants $`a_1`$ and $`a_2`$ and integers $`N_1`$ and $`N_2`$ characterizing the lattice sites. The components of $`𝐪`$ vary in the Brillouin zone (BZ) in the range $`[1/2,+1/2]`$. The dipole tensor (15) reads
$$𝐓(N_1,N_2)=\frac{1}{(N_1^2a_1^2+N_2^2a_2^2)^{5/2}}\left[\begin{array}{cc}(2N_1^2a_1^2N_2^2a_2^2)& 3N_1N_2a_1a_2\\ 3N_1N_2a_1a_2& (2N_2^2a_2^2N_1^2a_1^2)\end{array}\right]$$
(42)
Although for all figures the exact dynamical matrix is used, it is useful to consider the results with nearest neighbor (n.n.) lattice sums in (41) separately. This provides simple formulas for order– of– magnitude estimates.
$`C_{11}`$ $`=`$ $`\omega _1^2+2p_1[1cos(2\pi q_1)]p_2[1cos(2\pi q_2)]`$ (43)
$`C_{22}`$ $`=`$ $`\omega _2^2+2p_2[1cos(2\pi q_2)]p_1[1cos(2\pi q_1)]`$ (44)
$`C_{12}`$ $`=`$ $`\mathrm{\Omega }_{12}`$ (45)
where we introduced the interaction parameters $`p_i=\frac{2\beta N}{a_i^3}`$.
The convergence of the lattice sums $`S_{ik}`$ in the dynamical matrix is shown Fig.3. $`S_{ik}`$ is defined by
$$C_{ik}=\mathrm{\Omega }_{ik}+p_2S_{ik}$$
(46)
and depends only on $`𝐪`$ and the ratio $`a_1/a_2`$. Apart from the off-diagonal elements, which vanish in n.n. approximation, the error of the n.n. approximation is less than $`30\%`$.
### C Special Features of the Magneto– Phonon Spectrum
Fig.s 4-6 show the magneto– phonon The term magneto– phonon is attributed to the fact that the there is no exchange and there are harmonic forces between the oscillating individuals. One could also call them magneto– plasmons, if one wants to emphasize that it is only electrons which oscillate, and no nuclei. spectrum for circular dots on a rectangular lattice with $`a_1=2a_2`$. Because the two interaction parameters have a fixed ratio, it suffices to use one of them for characterizing the interaction strength. We chose the larger one $`p_2=p`$. For $`B=0`$ and isolated dots $`(p=0)`$, the two excitation modes are degenerate. If we tune up the interaction strength represented by $`p`$, a $`𝐪`$ dependent splitting appears (see Fig. 4).
This splitting is a manifestation of the dot interaction. For a certain critical $`p_{cr}`$ the lower mode becomes soft. This feature will be discussed below. There are points in the BZ, however, where the degeneracy for finite $`p`$ remains. These points will be investigated now. We demonstrated in Sect.II after formula (8) that necessary for degeneracy is $`C_{12}=0`$, i.e., the dynamical matrix must be diagonal. Then the points with degeneracy are defined by the condition $`C_{11}=C_{22}`$. As seen in (41), for circular dots $`\omega _1=\omega _2=\omega _0`$ this happens in the center of the BZ $`𝐪=0`$. The next question to be discussed is if there are other points with degeneracy. The first condition $`C_{12}=0`$, is fulfilled for all points on the surface of the BZ. The second condition must be investigated for special cases. We find, that for quadratic lattices $`a_1=a_2`$ with circular dots $`\omega _1=\omega _2`$ both modes are degenerate at the point $`𝐪=(1/2,1/2)`$. In the case shown in Fig.4 this point is somewhere between $`(1/2,1/2)`$ and $`(1/2,0)`$.
In n.n. approximation (44), however, this equation is even fulfilled on full curves in the BZ defined by $`p_1[1cos(2\pi q_1)]=p_2[1cos(2\pi q_2)]`$. In a cubic lattice, this is the straight lines $`q_2=\pm q_1`$. The contributions beyond n.n.s remove the exact degeneracy on this curve in the interior of the BZ, but leave a kind of anti-crossing behavior of the two branches.
An important parameter, which characterizes the influence of the dot interaction in circular dots, is the band width at $`B=0`$, i.e. the maximum splitting of the two branches due to dot interaction. (Remember that this splitting vanishes for noninteracting circular dots.) Assume $`a_1>a_2`$. Then the largest splitting for circular dots in n.n. approximation appears at $`𝐪=(0,1/2)`$ and has the amount
$$W=max(\mathrm{\Delta }E_+\mathrm{\Delta }E_{})=max(\omega _+\omega _{})=\sqrt{\omega _0^2+4p_2}\sqrt{\omega _0^22p_2}$$
(47)
For small dot interaction and in units $`\omega _0`$, this is proportional to the interaction parameter $`\frac{W}{\omega _0}3p_2`$.
We next discuss the appearance of soft modes. The question is, for which $`𝐪`$, $`B`$ and interaction parameter $`p`$ this happens. The general condition for vanishing of the lowest mode is $`C_{11}C_{22}=C_{12}^2`$ (see Sect. II). In this condition the magnetic field does not appear. For circular dots and with the definition (46) this equation reads
$$[\omega _0^2+p_2S_{11}][\omega _0^2+p_2S_{11}]=p_2^2S_{12}^2$$
(48)
After introducing a dimensionless critical interaction parameter $`P_{cr}=p_2/\omega _0^2`$, we obtain a quadratic equation for $`P_{cr}`$ which has the solution
$$P_{cr}=\frac{1}{2}\frac{tr}{det}\pm \sqrt{\left(\frac{1}{2}\frac{tr}{det}\right)^2\frac{1}{det}}$$
(49)
where $`det=S_{11}S_{22}S_{12}^2`$ and $`tr=S_{11}+S_{22}`$. For our numerical case $`a_1=2a_2`$ and n.n. interaction for $`S_{ik}`$ the lowest mode becomes soft at $`𝐪=(0,1/2)`$ and the critical interaction parameter is $`P_{cr}=1/2`$. Inclusion of lattice contributions beyond n.n. shifts $`P_{cr}`$ to 0.7543. The most important result of this paragraph is that lattice softening is independent of $`B`$ (see also Fig.s 5 and 6). The latter conclusion is exact within the range of validity of the Hamiltonian (31) and no consequence of any subsequent approximation or specialization.
## V Intra-dot– Excitations for N=2
Intra-dot excitations for circular dots in a cubic lattice and for $`N=2`$ can be calculated easily. We define the relative coordinate $`𝐫=𝐫_2𝐫_1`$, and assume that all dots are equivalent (also with respect to their environment). Then the indexes $`(n,\alpha )`$ can be chosen as $`(0,0)`$ and omitted. The relative Hamiltonian reads
$$H_{rel}=2\left\{\frac{1}{2m^{}}\left[𝐩+\frac{1}{2c}𝐀(𝐫)\right]^2+\frac{1}{2}𝐫𝐃𝐫+\frac{\beta }{2r}\right\}$$
(50)
where $`𝐩=i_𝐫`$ and
$$𝐃=\frac{1}{4}𝛀+\frac{\beta }{2}𝐓_0;𝐓_0=\underset{n,\alpha (0,0)}{}𝐓\left(𝐑_{n,\alpha }^{(0)}\right)$$
(51)
It is worth emphasizing that $`H_{rel}`$ contains a contribution from neighboring dots, originating from the interdot Coulomb interaction. A trivial angular dependent part can only be decoupled from $`H_{rel}`$, or, the 2-dimensional Schrödinger equation can be traced back to an ordinary radial Schrödinger equation, if the term $`𝐫𝐃𝐫`$ has the same circular symmetry as the intra-dot Coulomb term $`\beta /(2r)`$. Therefore we confine ourselves to circular dots on a cubic lattice, and we have
$$𝐓_0=\frac{1}{a^3}\underset{N_1,N_20,0}{}\frac{1}{(N_1^2+N_2^2)^{3/2}}𝐈\frac{4}{a^3}𝐈$$
(52)
where the simple result is in n.n. approximation. Using the interaction parameter $`p=2N\beta /a^3`$ (with $`N=2`$) defined above, we obtain
$$𝐃=\frac{1}{4}(\omega _0^2+2p)𝐈$$
(53)
In this way, dot interaction defines an effective confinement frequency $`\omega _{0,eff}^2=\omega _0^2+2p`$. This means that the c.m. excitations have to be calculated (or interpreted) with another confinement potential then the relative excitations. In our figures we present results for $`\omega _{0,eff}=0.2a.u.^{}`$, which agrees with the bare confinement potential used in Sect.IV and the mean value in Sect.III. Because our results are presented in units of $`\omega _0`$, they depend on $`\omega _0`$ only weakly through the differing influence of electron-electron interaction. For the absolute values, however, the influence of the dot interaction can be tremendous.
In the relative motion there is a coupling between orbital and spin parts through the Pauli principle. For $`N=2`$ and a circular effective confinement, Pauli principle demands that orbital states with even and odd relative angular momentum $`m_i`$ must be combined with singlet and triplet spin states, respectively ( see e.g. Ref.). For the c.m. motion there is no interrelation between orbital and spin part because the c.m. coordinate is fully symmetric with respect to particle exchange. Consequently, any c.m. wave function can be combined with a given spin eigen function. The only spin dependent term in the total energy considered here is the Zeeman term, which reads in our units
$$\frac{E_B}{\omega _0}=0.913410^2g_s\frac{B[Tesla]}{\omega _0[a.u.^{}]}\frac{M_s}{2}$$
(54)
where we used $`g_s=0.44`$ for the gyro-magnetic factor of GaAs from Ref.. The total spin quantum number is $`M_s=0`$ for the singlet state and $`0,\pm 1`$ for the triplet state. One of the most interesting points in quantum dot physics is that the total orbital angular momentum of the ground state depends on the magnetic field (see e.g. Ref.s,). This feature is a consequence of electron– electron interaction. For our parameter values, the relative orbital angular momentum of the ground state $`m_i`$ changes from 0 to –1, from –1 to –2, and from –2 to –3 at $`B=1.250`$, $`4.018`$, and $`5.005Tesla`$. This corresponds to a sequence $`M_s`$=0,+1,0,+1 for the spin quantum number. Figures 7a-c show the excitation frequencies for three B-values lying within the first three regions. $`m_f`$ is the relative orbital angular momentum of the final state. All excitations are included irrespective of selection rules. For dipole transitions only two of them would remain (the lowest excitation with $`m_f=m_i\pm 1`$). For $`B=0`$, the lowest excitation energy (in units $`\omega _0`$) for noninteracting electrons would be 1. As seen in in Fig.7a, electron– electron interaction decreases this value by at least a factor of $`1/2`$. The same holds qualitatively for finite $`B`$. This is connected to the fact, that the ground state depends on $`B`$. Let’s consider an example. For $`B=1.250Tesla`$ the ground state switches from $`m_i=0`$ to $`m_i=1`$. This implies that for $`B`$ approaching this transition field from below, the excitation energy for dipole allowed transition from $`m_i=0`$ to $`m_f=1`$ converges to $`0`$. In other words, there is a level crossing at the the transition field. Therefore, very small transition energies and switching of the ground state are connected.
For a qualitative understanding, Figures 7a-c can be used together with Figures 1a, 1b, 4, and 5 to investigate the relative position of collective and intra-dot excitations. The conclusion is that for small dot interaction (for $`p`$ well below $`p_{cr}`$), the lowest intra-dot excitation energies lie well below the lowest c.m. excitations. Apart from using a different terminology, this conclusion agrees with the experimental findings in Ref..
Fig.s 7b and c, which belong to finite $`B`$, show the Zeeman splitting. All transition energies to final states with odd $`m_f`$ are triplets because the corresponding spin state is a triplet state. The thin lines of a triplet belong to spin– flip transitions.
In Fig.8 the $`B`$– dependence of the lowest excitation energies is shown. It is clearly seen that the curves exhibit a kink at those $`B`$– values, where the ground state configuration changes. The size of the kink decreases with increasing $`B`$. If this kink could be resolved experimentally (e.g. by electronic Raman spectroscopy), it would be a direct indication for the change of the ground state configuration, and thus an experimentally observable consequence of electron– electron interaction.
## VI Summary
We solved the Schrödinger equation for a lattice of identical parabolic (but not necessarily circular) quantum dots with Coulomb interaction (in dipole approximation) between the dots. We provide an overview over the state of art of these systems which includes the results of former publications. References can be found in the text.
* Similar to single dots, the center of mass coordinates of all dots can be separated from the relative coordinates. Only the c.m. coordinates of different dots are coupled to each other. The relative coordinates of different dots are neither coupled to each other nor to the c.m. coordinates.
* This gives rise to two types of excitations: two collective c.m. modes per dot and and a complex spectrum of intra-dot excitations. In periodic arrays only the collective c.m. modes show dispersion. Intra-dot excitations are dispersion-less.
* The c.m. system can be solved exactly and analytically providing magneto– phonon excitations characterized by a certain wave number $`𝐪`$ within the Brillouin zone. For $`𝐪=0`$ and one dot per unit cell, interdot interaction does not have any influence on the c.m. excitations.
* All dipole allowed excitations (seen in FIR experiments) are not influenced by the dot interaction.
* Interdot interaction between two dots influences the spectrum through a single parameter $`p=2N\beta /a^3`$, where $`a`$ is the distance between the dots, $`N`$ the number of electrons per dot and $`\beta `$ the inverse background dielectric constant.
* If $`p`$ exceeds a certain critical value $`p_{cr}`$, the lowest c.m. mode becomes soft leading to an instability. This transition is independent of the magnetic field.
* For B=0 and and one circular dot per unit cell, the two c.m. modes are not only degenerate in the middle of the Brillouin zone, but also at some points on the surface. If we use the n.n. approximation for the lattice sums in the dynamical matrix, degeneracy is maintained even on full curves in the Brillouin zone.
* Intra-dot excitations have to be calculated from an effective confinement. In circular dots with a cubic environment in nearest neighbor approximation the effective confinement frequency reads $`\omega _{0,eff}^2=\omega _0^2+2p`$. This effective confinement differs from that for the c.m. motion.
* For $`p`$ well below $`p_{cr}`$, the lowest intra-dot excitations are much smaller than the lowest collective excitations.
* The intra-dot excitation energies versus magnetic field exhibit kinks at those fields, where the angular momentum of the ground state changes.
In the Appendix we prove a Kohn Theorem for dot arrays with Coulomb interaction between the dots without the dipole approximation. The individual confinement potentials can be arbitrarily arranged and can carry different electron numbers, but have to be described by identical confinement tensors. This means that for breaking Kohn’s Theorem in dot arrays, we have to have at least two different confinement species.
## Appendix
We are going to prove that for an arbitrary array <sup>§</sup><sup>§</sup>§The dot centers can be arranged arbitrarily. of identical parabolic quantum dot potentials The confinement tensors $`𝛀`$ of all dots must be equal. in an homogeneous magnetic field: i) the total c.m. degree of freedom can be separated from the rest, ii) the total c.m. Hamiltonian is not influenced by Coulomb interaction, and iii) the eigenvalues of the total c.m. Hamiltonian are independent of the electron number $`N`$ in each dot The electron number in different dots can be different. .
The Hamiltonian $`H=H^{(0)}+V`$ consists of an one–particle term $`H=H^{(0)}`$ and the Coulomb interaction between all electrons $`V`$. The dot centers are located at $`𝐑_\alpha ^{(0)}`$ and the electron coordinates are denoted by $`𝐫_{i\alpha }=𝐑_\alpha ^{(0)}+𝐮_{i\alpha }`$. Then we have
$$H^{(0)}=\underset{i\alpha }{}\left\{\frac{1}{2m^{}}\left[𝐩_{i\alpha }+\frac{1}{c}𝐀(𝐑_\alpha ^{(0)}+𝐮_{i\alpha })\right]^2+\frac{1}{2}𝐮_{i\alpha }𝐂𝐮_{i\alpha }\right\}$$
(55)
First of all we shift the gauge center for each electron into the middle of the corresponding dot using an unitary transformation similar to (16). This transforms the shift $`𝐑_\alpha ^{(0)}`$ in the argument of the vector potential away. Next we drop the index $`\alpha `$ in (55) so that the index ’$`i`$’ runs over all electrons in all dots. Now we perform a transformation to new coordinates $`\stackrel{~}{𝐮}_i`$
$$𝐮_i=\underset{k}{}Q_{ik}(\sqrt{N}\stackrel{~}{𝐮}_k);(\sqrt{N}\stackrel{~}{𝐮}_i)=\underset{k}{}Q_{ki}^{}𝐮_k$$
(56)
where $`Q_{ik}`$ is an unitary matrix. This implies
$$𝐩_i=\underset{k}{}Q_{ik}^{}(\frac{\stackrel{~}{𝐩}_k}{\sqrt{N}});(\frac{\stackrel{~}{𝐩}_i}{\sqrt{N}})=\underset{k}{}Q_{ki}𝐩_k$$
(57)
It is possible to choose for the first column $`Q_{k1}=\frac{1}{\sqrt{N}}`$. The other columns need not be specified. Then $`\stackrel{~}{𝐮}_1=(1/N)_i𝐮_i=𝐔`$ is the c.m. of all elongations, or, the c.m. of the electron coordinates with respect to the weighted center of the dot locations $`𝐑^{(0)}=(1/N)_\alpha N_\alpha 𝐑_\alpha ^{(0)}`$, where $`N_\alpha `$ is the number of electrons in dot $`\alpha `$. The corresponding canonical momentum $`\stackrel{~}{𝐩}_1=(1/i)_{\stackrel{~}{𝐮}_1}=𝐏`$ is the c.m. momentum. Inserting our transformation into (55) provides
$$H^{(0)}=\underset{i}{}\left\{\frac{1}{2m^{}}\left[\frac{1}{\sqrt{N}}\stackrel{~}{𝐩}_i+\frac{\sqrt{N}}{c}𝐀(\stackrel{~}{𝐮}_i)\right]^2+\frac{N}{2}\stackrel{~}{𝐮}_i𝐂\stackrel{~}{𝐮}_i\right\}$$
(58)
The term $`i=1`$ in (58) is the (separated) c.m. Hamiltonian
$$H_{c.m.}=\frac{1}{2m^{}}\left[\frac{1}{\sqrt{N}}\stackrel{~}{𝐏}+\frac{\sqrt{N}}{c}𝐀(\stackrel{~}{𝐔})\right]^2+\frac{N}{2}\stackrel{~}{𝐔}𝐂\stackrel{~}{𝐔}$$
(59)
which agrees with (10). Clearly, the Coulomb interaction $`V`$ in $`H`$ is independent of the c.m., and does not contribute to $`H_{c.m.}`$. For the independence of the eigenvalues of $`N`$ see the discussion following (10).
This prove, in particular the step from (55) to (58), is not correct if the dot confinement tensor $`𝐂`$ depends on $`\alpha `$ (or ’$`i`$’ in the changed notation). Therefore, all dots must have the same $`𝐂`$, but can have different electron numbers $`N_\alpha `$. In other words, the total c.m. excitations in dot arrays, which are seen in FIR spectra, are not affected by the e– e– interaction, if and only if all confinement tensors $`𝐂`$ are equal. On the other hand, if we want to observe e– e– interaction in the FIR spectra and break Kohn’s Theorem, we have to use dot lattices with at least two different confinement tensors. The simplest way to implement this is using a lattice with two non-circular dots per cell, which are equal in shape, but rotated relative to each other by 90 degrees.
## VII Acknowledgment
I am indebted to D.Heitmann, H.Eschrig, and E.Zaremba and their groups for very helpful discussion and the Deutsche Forschungs– Gemeinschaft for financial funding. |
warning/0002/astro-ph0002497.html | ar5iv | text | # Limits on Hot Intracluster Gas Contributions to the Tenerife Temperature Anisotropy Map
## 1 Introduction
The discovery of Cosmic Microwave Background (CMB) temperature anisotropies (Smoot et al. 1992; Wright et al. 1992) immediately prompted the question of their origin. In most cosmological scenarios temperature anisotropies are expected to arise with the growth of matter density perturbations (Peebles 1980); but it was also suggested (Hogan 1992) that the anisotropies could have been originated by inverse Compton scattering from hot diffuse clouds of electrons in nearby superclusters (see Birkinshaw 1999 for a review on the Sunyaev-Zel’dovich effect). Right after Hogan’s suggestion, Boughn & Jahoda (1993) and Bennett et al. (1993) searched for non cosmological signal contributions to the COBE/DMR 2 year sky maps. If the Sunyaev-Zel’dovich (SZ) component is significant at some level and the distribution of hot electrons is traced by rich clusters or any other extragalactic source, then temperature anisotropies should be correlated with maps constructed from extragalactic source surveys. Banday et al. (1996) performed this analysis for the COBE/DMR 4 year map. Their study did not find a statistically significant contribution, reassuring the idea that temperature anisotropies were of cosmological origin. Similar cross-correlation analysis have been carried out by several authors to estimate the rms level of the Galactic emission present in different CMB maps (Kogut et al. 1996; Kneissl et al. 1997; de Oliveira-Costa et al. 1997, 1998, 1999).
In this article we study the contribution of nearby structures on the temperature anisotropies measured by the Tenerife map (Gutiérrez et al. 2000). Two effects associated with the local matter distribution can give rise to anisotropies in the CMB: the non-linear evolution of gravitational structures (Rees & Sciama 1968) and the mentioned SZ effect. The latter is expected to give the largest contribution, and we shall center our analysis on it. Since clusters are known to posses ionized gas at temperatures about $`10^8`$K they are expected to give the largest contribution to the SZ effect. Cocoons of radio galaxies (Yamada, Sugiyama & Silk 1999) or other foreground sources subtend a very small angle and, due to beam dilution, have a very small contribution at the scales probed by the Tenerife experiment. Therefore, we shall look for correlations between the Tenerife CMB data and template maps generated from cluster catalogues; in particular, we use the ACO (Abell, Corwin & Olowin 1989), and ROSAT PSPC catalogues (Vikhlinin et al. 1998). To test the hypothesis that the CMB signal is generated by hot diffuse gas distributed on supercluster scales we also include a supercluster catalogue (Einasto et al. 1994) in our analysis. Finally, we also constructed template maps from the HEAO 1 A-1 map of X-ray sources (Kowalski et al. 1984) that should trace the distribution of hot electrons. In Sec. 2 we describe the statistical methods used for comparing data on temperature anisotropies with template maps constructed from different surveys. In Sec. 3 we describe the catalogues used and how the template maps were elaborated. Finally, in Sec. 4 we present and discuss our main results.
## 2 Statistical Method.
At any given frequency, the CMB anisotropy map can be considered a superposition of contributions of cosmological origin, $`T_{CMB}`$, astrophysical origin, $`\alpha M`$, and instrument noise, $`N`$: $`T=T_{CMB}+\alpha M+N`$. $`\alpha `$ is a conversion factor -to be determined- that gives the amplitude of the contribution of foreground sources to the CMB temperature anisotropies. As mentioned in the introduction, the spatial distribution of clusters should trace that of the electrons. From cluster catalogues we shall construct template maps of the temperature anisotropies induced on the CMB spectrum by the hot gas traced by clusters. We shall name template map the term $`\alpha M`$. Let us remark that we do not have a map of the distribution of gas on the nearby Universe. Our hypothesis that clusters trace the gas distribution means that we expect the autocorrelation function of the template map to be rather similar to that of the hot gas even though we ignore the exact gas distribution. Therefore, and unlike de Oliveira-Costa et al. (1999), our analysis shall be based on comparing correlation functions and not the maps themselves.
Assuming that the contribution of foreground sources is uncorrelated with the cosmological signal and noise in the temperature anisotropy map, then the cross correlation of the CMB and the template maps $`C_{TM}(\theta )`$ is related with the template map autocorrelation function $`C_{MM}(\theta )`$ as: $`C_{TM}(\theta )=\alpha C_{MM}(\theta )`$. A best-fit value of $`\alpha `$ is obtained by minimizing (Banday et al. 1996)
$$\chi ^2=\underset{ij}{}[C_{TM}(\theta _i)\alpha C_{MM}(\theta _i)]M_{ij}^1[C_{TM}(\theta _j)\alpha C_{MM}(\theta _j)].$$
(1)
In this expression $`M_{ij}`$ is the covariance matrix of the cross-correlation functions (Ganga et al. 1993) defined as follows: $`M_{ij}=<[C(\theta _i)<C(\theta _i)>][C(\theta _j)<C(\theta _j)>]>`$, with $`\theta _i,\theta _j`$ two arbitrary angular separations in the sky. $`C(\theta )`$ is the cross correlation of the template map and one single realization of the observed sky. Realizations of the sky were performed in two different ways: (a) at each measured temperature we add a random realization of a gaussian distributed noise with zero mean and the variance at that point. (b) We performed Monte Carlo CMB simulations of the Tenerife data drawn from a gaussian distribution with variance $`C_l=6C_2/l(l+1)`$ at each multipole, normalized to $`Q_{rmsPS}=20\mu `$K. We assumed a Harrison-Zel’dovich power spectrum for the primordial fluctuations since, together with the previous normalization, is a good approximation at the scales probed by the Tenerife experiment (Gutiérrez et al. 2000). To each point in the CMB map we add a realization of the noise as in (a). In both cases, the average $`<..>`$ was obtained from a thousand realizations.
The minimum-variance estimate is:
$$\widehat{\alpha }=\underset{ij}{}\frac{C_{MM}(\theta _i)M_{ij}^1C_{TM}(\theta _j)}{_{ij}C_{MM}(\theta _i)M_{ij}^1C_{MM}(\theta _j)}$$
(2)
with formal error
$$\sigma _{\widehat{\alpha }}=(\underset{ij}{}C_{MM}(\theta _i)M_{ij}^1C_{MM}(\theta _j))^{1/2}.$$
(3)
The approach (a) described above does not include sample variance. We estimated the associated error bar by performing a thousand Monte Carlo realizations of the CMB sky and finding $`\alpha `$ from the correlation with the template maps. As expected, the average value of $`\alpha `$ was zero. The dispersion around this mean, $`\sigma _s`$, is a measure of both cosmic variance and the variance coming from random alignments. On the other hand, the approach (b) includes all contributions to the variance in the estimate of $`\widehat{\alpha }`$.
## 3 Data and Template Maps.
The results of the Tenerife CMB experiments are presented in Gutiérrez et al. (2000). The observations were performed in two frequencies: 10 and 15 GHz covering 5000 and 6500 square degrees, respectively. The experiments are sensitive to multipoles $`l=1030`$ which corresponds to the Sachs-Wolfe plateau of the radiation power spectrum. The experiment measures strips in right ascension separated by $`2.5^o`$ in declination. The 15GHz map is made of 8 strips that spans a region on the sky from 8h to 18h in R.A. and from $`27.5^o`$ to $`45^o`$ in Dec. The 10GHz is slightly smaller with only 5 strips running from $`32.5^o`$ up to $`42.5^o`$. The map is in the North Galactic hemisphere and has a galactic latitude $`b20^o`$. The experiment uses a double-differencing technique to measure, with a $`5^o`$ FWHM beam, points separated $`8.1^o`$ in R.A. For 15GHz, the band power of the CMB signal is $`\mathrm{\Delta }T_l=30_{11}^{+15}\mu `$K, including a possible contaminating effect due to the diffuse Galactic component. The r.m.s. temperature anisotropy at $`5^o`$ is $`\sigma _{TEN,10GHz}=43\mu `$K. At 15GHz, $`\sigma _{TEN,15GHz}=32\mu `$K. The sensitivity at 10 and 15 GHz was $`31\mu `$K and $`12\mu `$K, respectively, in a beam-size region.
As explained in the previous section, we assume that clusters trace the spatial distribution of the hot gas. Cluster surveys select members according to a given criteria. Therefore, different catalogues have different selection biases. For each catalogue we shall elaborate a template map to compare with the Tenerife CMB data. Let us briefly describe the ones that will be used. The ACO all-sky catalogue contains 4073 rich clusters of galaxies, each having at least 30 members within magnitude range $`m_3`$ to $`m_3+2`$ ($`m_3`$ is the magnitude of the third brightest cluster member) and each with redshift less than 0.2. The HEAO 1 A-1 catalogue is essentially a catalogue of ACO clusters with X-ray emission in the energy range $`0.520`$keV. For several nearby clusters the SZ effect has been measured (Birkinshaw 1998). Therefore, this catalogue traces the extragalactic objects known to be sources of SZ. In this respect, it will be interesting to compare the results obtained by cross-correlating each of these templates with the Tenerife map. By including clusters that do not contribute significantly to the SZ effect, we could have diluted the signature of the hot gas in the cross-correlation between the ACO template and Tenerife.
ROSAT is a catalogue of X-ray selected objects. It includes from poor groups till rich clusters of galaxies. These clusters were serendipitously detected in the ROSAT PSPC high Galactic latitude pointed observations ($`b30^o`$). The satellite covers a large energy range ($`0.12`$keV) in the soft X-ray band. The cluster redshifts range from $`z=0.015`$ to $`z>0.5`$ in the area of the sky covered by Tenerife. The HEAO 1 A-1 and ROSAT catalogues are less sensitive than optical catalogues to projection effects and could detect ”failed clusters” were galaxy formation was suppressed.
The Tenerife experiment operates on the Rayleigh-Jeans regime and the effect of the hot electrons is to produce a decrement on the radiation temperature. Therefore, if the experiment has detected any contribution from hot gas, cold spots in the data and the template maps should be correlated. The anisotropy depends linearly on the central electron temperature, cluster core radius and electron density: $`\delta T/T_or_cT_en_e`$ (Zel’dovich & Sunyaev 1969). The exact relation depends on the cluster density profile but this is of no significance since clusters are unresolved by the antenna. The parameters $`r_c,n_e`$ and $`T_e`$ scale with the cluster mass (Bower 1997). To elaborate a template map we assume that not only clusters trace the gas distribution but also that the cluster richness is a measure of the cluster mass, and, consequently, of its size, gas content and electron temperature. We only had information on the richness of ACO and HEAO 1 A-1 clusters. For them, we constructed a template map by assigning a number to each pixel: zero if there was no cluster, and a contribution proportional to the richness $`(Richness)^n`$ if there was a cluster. Since we did not know how to scale the cluster size, gas density and electron temperature with richness class we tried exponentiating the richness class to three different powers: $`n=0,1,2`$. Finally, this pixel map was convolved with the Tenerife window function. To avoid boundary effects, we included objects within $`15^o`$ of the region probed by Tenerife. Like in Bennett et al. (1993) we found no significative effect: while for the ACO clusters larger $`n`$ led to diminishing the cross-correlation and consequently the value of $`\alpha `$, for the HEAO 1 A-1 catalogue, the opposite effect was observed. However, in all cases the effect was minute and well within the error bars. We shall quote our results for $`n=0`$, when all clusters contribute equally to the SZ effect. The template map of the ROSAT catalogue was constructed in a similar manner but without scaling with richness class.
Finally, we also included in our analysis the supercluster catalogue of Einasto et al. (1994). This catalogue was elaborated from the distribution of rich clusters of galaxies up to redshift $`z=0.1`$, extracted from the ACO catalogue described above. For each supercluster, the number of cluster members, center position, average distance $`D`$, extent in supergalactic coordinates and length $`L`$ in Mpc are given. In our analysis, twenty two superclusters with typical angular sizes between 5 and $`10^o`$ were included. Contrary to clusters, superclusters can not be considered point-like. Therefore, different hypothesis about the gas distribution could lead to different results. As a first approximation, we took the gas distributed homogeneously on a sphere of size $`2\mathrm{tan}^1(L/2D)`$. This template is very convenient in order to check Hogan (1992) hypothesis about the local origin of temperature anisotropies. We called this template ”superclusters with homogeneous gas distribution”. We checked that the correlation level did not depend on the gas distribution by chosing a model with a density profile:
$$n(r)=\frac{n_e}{1+(r/r_c)^2}$$
(4)
where $`r_cL/10`$ is a fiducial radius. We called this template ”superclusters with concentrated gas”. Finally, both templates were convolved with the Tenerife beam pattern before performing the correlation analysis.
## 4 Numerical Results and Discussion.
Table 1 summarizes the results of our analysis. After substracting the mean and normalizing the templates to unit variance, the autocorrelations and cross correlations were computed given equal weight to each pixel. We tried different angular bins, from 1<sup>o</sup> to 5<sup>o</sup>, and computed the correlation function out to $`20^o`$ and $`30^o`$. No significant differences were found. In Table 1, the results are quoted for correlation functions in bins of $`3^o`$ out to $`21^o`$.
The SZ effect will generate approximately equal and negative contributions at 10 and 15GHz. If the signal detected is real one should expect equal and positive values of $`\alpha `$ at the two frequencies. In Table 1 we give $`\widehat{\alpha }`$, $`\sigma _{\widehat{\alpha }}`$ as given by eqs. (2) and (3), and $`\chi ^2`$ per degree of freedom (dof) in the two approaches described in Sec. 2. For simplicity, we termed (a) “without sample variance” and (b) “with sample variance”. In the case (a), we also give the error associated with sampling variance $`\sigma _s`$, which should be added in quadrature with $`\sigma _{\widehat{\alpha }}`$. As the maps were normalized to unit variance, $`\alpha \sigma _{TEN}`$ gives the SZ component of the CMB map in thermodynamic units. For each template we calculate the cross-correlation with the 15GHz and 10GHz maps. Since the latter covers a smaller fraction of the sky, we also correlate a reduced 15GHz map (denoted by 15c in the table) cut to the size of the 10GHz map to eliminate the bias introduced by the different sky coverage.
No significative detections (larger than $`2\sigma `$) were found by either of the two methods. We always found negative values of $`\alpha `$, i.e., cold spots in the template map correlate with hot spots in the data, contrary to what one would expect if there was a significant SZ contribution. The largest signal was obtained at the 15c ROSAT template map. Since the amplitude of the SZ effect does not change much at the Tenerife frequencies, consistency would require a fluctuation of the same order to be present at 10GHz. Furthermore, when sample variance was not included in the covariance matrix, the best-fit was never a good fit. Only when it was included $`\chi ^2/dof`$ became of order unity. The low quality of the fit can be understood by looking at Figure 1, where we plot the autocorrelation of the template maps (Fig 1a) and their cross correlation with the data on 15GHz (Fig 1b). The dashed line corresponds to the autocorrelation and cross correlation of the ACO catalogue, the long dashed line to the HEAO 1 A-1 catalogue, the dot-dashed line to the ROSAT catalogue. Thick and thin solid lines correspond to superclusters with uniform and concentrated gas, respectively. While the autocorrelation functions are rather similar in all cases, the cross-correlation differ substantially in shape. When the sample variance is not included in the covariance matrix, $`\sigma _{\widehat{\alpha }}`$ is small and the difference in shape implies a large $`\chi ^2/dof`$. Only when sample variances are included, the error bars are much larger and the fit to the data improves. To conclude, we can only set upper limits on the value of $`\alpha `$. Taking the results on 15GHz, we limit $`\alpha 0.24`$ at the 99% confidence level.
For an experiment with such a beam width like Tenerife, one could not expect to find a large correlation between data and templates. For example, some clusters in the ACO catalogue have been found to produce temperature fluctuations of the order $`100\mu `$K (Birkinshaw 1999). But they subtend an angular scale of few arcminutes and as a result the SZ signal is diluted by the large solid angle covered by the Tenerife beam. Still, the Tenerife data limits the contribution of nearby clusters and superclusters to be smaller than $`8\mu `$K at 99% confidence level. The mean Comptonization parameter at $`y=\frac{\mathrm{\Delta }T}{2T_o}1.5\times 10^6`$ at the same level of confidence. Our results are slightly more restrictive than those previously found by Banday et al. (1996). Let us remark that $`y`$ obtained above only limits the contribution due to nearby superclusters, while the COBE result of Mather et al. (1994) applies to the contribution of all structures located between the last scattering surface and the observer. To conclude, this study, like previous ones based on the COBE/DMR data, indicate that most of the signal measured by Tenerife is not of extragalactic origin but cosmological.
We thank R. Rebolo for many useful discussions and comments. F.A.B acknowledge the financial support of the University of La Laguna - Banco de Santander. F.A.B. and C.H.M. acknowledge the hospitality of the I.A.C. where most of this work was carried out. |
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