id
stringlengths 27
33
| source
stringclasses 1
value | format
stringclasses 1
value | text
stringlengths 13
1.81M
|
|---|---|---|---|
warning/0506/math0506180.html | ar5iv | text | # Constructions in public-key cryptography over matrix groups
## Introduction
One of the oldest cryptographical problems consists in constructing of a key agreement protocol. Roughly speaking it is a multi-party algorithm, defined by a sequence of steps, specifying the actions of two or more parties in order a shared secret becomes available to two or more parties. Probably the first such procedure based on abelian groups is due to Diffie-Hellman one (see ). In fact, it concerns automorphisms of abelian (even cyclic) groups induced by taking to a power. Some generalizations of this protocol to non-abelian groups (in particular, the matrix groups over some rings) were suggested in where security was based on an analog of the discrete logarithm problems in groups of inner automorphisms. Certain variations of the Diffie-Hellman systems over the braid groups were described in ; there several trapdoor one-way functions connected with the conjugacy and the taking root problems in the braid groups were proposed. Recently, a general scheme for constructing key agreement protocols based on algebraic structures was proposed in . In principle, it enables us to construct such protocols for non-abelian groups and their automorphisms induced by conjugations. In this paper we generalize to the non-abelian case the Diffie-Hellman protocol, construct multi party procedure for the protocol , and analyze the security of both protocols realized in matrix groups over rings.
The question on finding probabilistic public-key cryptosystems in which the decryption function has a homomorphic property goes back to (see also ). In such a cryptosystem the spaces of messages and of ciphertexts are algebraic structures $`G`$ and $`H`$ and the decryption function $`D:GH`$ is a homomorphism. A number of such cryptosystems is known for abelian groups, e.g. the quadratic residue cryptosystem and its generalization for highest residues (see also an overview in ). In most of them the security is based on the intractability of theoretical number problems close to the integer factoring. Recently, several homomorphic cryptosystems were constructed for infinite (but finitely presented) groups, see and references there. In this paper we construct one more homomorphic cryptosystem with $`G`$ being a free group the trapdoor of which uses a secret permutation of the generators of $`G`$.
The third problem considered in this paper is how to produce instances for cryptosystems based on computations with matrix groups over rings. In contrast to numerous theoretical cryptosystems where there is a lot of efficient algorithms to generate integers with given properties (e.g., the pairs of two distinct large primes of the same bit size used in the quadratic residue cryptosystem), it is not clear a priory how to find efficiently matrix groups in which some problems (like membership or conjugacy) arising in cryptography are computationally difficult. We propose a general scheme for solving this problem and give a specialization of this scheme for matrix groups over finite commutative rings.
In Section 1 we study key agreement protocols between two parties (named usually Alice and Bob). The security of the Diffie-Hellman protocol relies on the difficulty of the following transporter problem: having an action $`G\times VV`$ of a group $`G`$ on a set $`V`$ for given $`u,vV`$ to find $`gG`$ (provided that it does exist) such that $`(g,u)v`$. In case of $`V`$ being a cyclic group of order $`n`$ and $`G`$ being a group acting on $`V`$ by taking a power one arrives to the discrete logarithm problem (usually, $`n`$ is taken to be prime). The security of the key agreement protocol of (see also Subsection 1.1) relies on the difficulty of the conjugacy problem with respect to a subgroup of $`G`$. In Subsection 1.1 we extend the construction of to multi-party key agreement protocol. Then in Subsection 1.2 we design another generalization of the Diffie-Hellman protocol to actions of groups $`G`$ which satisfy a certain identity. Clearly, any abelian group satisfies the identity $`aba^1b^1=1`$ and more generally, any solvable group with a fixed length of its derived series satisfies an appropriate commutator identity. The security of our protocol again relies on the difficulty of the transporter problem for a suitable action of $`G`$.
In Section 2 we consider homomorphic public key cryptosystems (see e.g. ) in which the decrypting function (known to Alice) is a group homomorphism $`f:GH`$ where the groups $`H,G`$ play the roles of the spaces of plain and ciphertext messages respectively. Usually, the security of a homomorphic cryptosystem relies on the difficulty of the problem of the membership to a normal subgroup of $`G`$ (here, the kernel of $`f`$). Also in Section 2 we describe a homomorphic cryptosystem in which as $`G`$ a free group is taken. This cryptosystem modifies one from where as $`G`$ a subgroup of the modular group $`SL_2()`$ was considered. The security of this cryptosystem relies on the difficulty of a certain word problem. A private key of Alice is an appropriate permutation of the generators of the free group $`G`$, this differs our cryptosystem from the one produced in .
The crucial role in the classical cryptographic constructions (like RSA, discrete logarithm or quadratic residue ) plays the natural action of the group $`Aut(_n^{})`$ on the group $`_n^{}`$. So, varying $`n`$ one gets a mass pool of instances for cryptographic primitives. This action is a special case of the natural action of the group $`Aut_R(V)`$ (viewed as a matrix group) on the free module $`V`$ over the ring $`R`$. In this paper we propose a construction of a pool of matrix groups instances for cryptographic primitives (Subsection 3.2). The security of these instances relies on the difficulty of certain problems on matrix groups (e.g. the membership to a subgroup or the conjugacy with respect to a subgroup). For the complexity of such problems few results were established in case of matrix groups over fields ; for matrix groups over arbitrary rings much less is known.
The common way in cryptography of producing a trapdoor and a cryptosystem, is to generate a private key departing from a pair of primes $`p,q`$, while their product $`n=pq`$ plays the role of a public key. In our scheme (see Subsection 3.1) as a private key we take a rooted tree whose leaves being furnished with specially chosen (non-abelian, in general) groups. We assume that Alice has in possession such representations of these groups which allow her to solve efficiently a problem lying in the background of a cryptosystem (like membership or conjugacy). Internal vertices of the tree are endowed with certain operations on groups which allow one to assign recursively a group to each vertex of the tree starting with its leaves. At the end of the recursion a group is assigned to the root, and this group plays the role of a public key. This scheme is also modified to produce a homomorphism of matrix groups as a public key. In Subsection 3.2 we give a realization of this general scheme in finite matrix groups.
The similarity of the common constructions in cryptography based on commutative groups (say. $`_n^{}`$) with our construction (relying on finite matrix groups) allows us to call the latter type of constructions the non-commutative cryptography.
## 1 Group theoretical key agreement protocol
### 1.1 A multi-party protocol.
The following group theoretical variant of key agreement two-party protocol was proposed in . Let $`G`$ be a group, and to two parties $`A`$ and $`B`$ are assigned their subgroups
$$G_A=a_1,\mathrm{},a_m,G_B=b_1,\mathrm{},b_n.$$
(1)
The group $`G`$ and the elements $`a_i`$, $`b_j`$ are publicly known. The parties $`A`$ and $`B`$ choose secret elements $`aG_A`$ and $`bG_B`$ and transmit to each other the collections
$$X_B=\{a^1b_ja\}_{j=1}^n,X_A=\{b^1a_ib\}_{i=1}^m$$
respectively. Since $`A`$ (resp. $`B`$) has a representation of the element $`a`$ (resp. $`b`$) via generators $`a_1,\mathrm{},a_m`$ (resp. $`b_1,\mathrm{},b_n`$), then $`A`$ (resp. $`B`$) can compute a representation of the element $`b^1ab`$ (resp. $`a^1ba`$) via elements of the set $`X_A`$ (resp. $`X_B`$). Thus $`A`$ and $`B`$ have a common key
$$a^1(b^1ab)=[a,b]=(a^1ba)^1b.$$
An obvious necessary condition for this protocol to be secure is that the set of all such commutators with $`aG_A`$, and $`bG_B`$ would contain at least two elements.
Let us describe a generalization of the group theoretical key agreement protocol for $`s`$ parties with $`s2`$ and a single public communicating channel. Without loss of generality we assume that $`s=2^t`$ for some $`t1`$, for otherwise in the recursive construction below we divide the parties into two unequal subsets which leads just to slight changing the notation. As in the case $`s=2`$ the groups $`G_1,\mathrm{},G_sG`$ of the parties are given publically by their sets of generators. At the initial step the $`i`$th party chooses a secret key $`a_iG_i`$, $`i=1,\mathrm{},s`$. Let $`S_1`$ and $`S_2`$ be disjoint $`s/2`$-subsets of the set $`\{1,\mathrm{},s\}`$. Then given $`u=1,2`$ the parties from $`S_u`$ recursively construct the common key $`K_uG`$, such that for all $`iS_u`$ there exist integers $`\epsilon _{i,j}\{1,+1\}`$ and $`1m_is/2`$, and certain elements $`B_{i,1},\mathrm{},B_{i,m_i}\{a_j:jS_{u,i}\}`$ with $`S_{u,i}=S_u\{i\}`$, for which we have
$$K_u=(B_{i,1}^1a_i^{\epsilon _{i,1}}B_{i,1})\mathrm{}(B_{i,m_i}^1a_i^{\epsilon _{i,m_i}}B_{i,m_i}).$$
By recursion we can assume that the $`i`$th party knows the elements $`B_{i,j}^1aB_{i,j}`$ for all $`j`$ and for all chosen generators $`a`$ of the group $`G_i`$ (and thereby, it knows $`B_{i,j}^1a_iB_{i,j}`$), but does not necessary know $`B_{i,j}`$. At this point the party $`iS_u`$ sends the elements $`B_{i,j}^1aB_{i,j}`$ for all the chosen generators $`a`$ of the group $`G_i`$ to a certain party from the set $`S_u^{}`$ with $`u^{}=3u`$ and asks for the elements $`K_u^{}^1B_{i,j}^1aB_{i,j}K_u^{}`$. Then for $`u=1`$ the $`i`$th party computes the element
$$[K_1,K_2]=K_1^1(K_2^1K_1K_2)=K_1^1(K_2^1(B_{i,_{m_i}}^1a_i^{\epsilon _{i,m_i}}B_{i,m_i})K_2)\mathrm{}(K_2^1(B_{i,1}^1a_i^{\epsilon _{i,1}}B_{i,1})K_2).$$
Similarly, for $`u=2`$ the $`i`$th party computes the element $`[K_1,K_2]=(K_1^1K_2K_1)^1K_2`$. Thus this element can be chosen as the common key for all parties. It is easy to see that the $`i`$th party computes the common key in $`O(s|a_i|)`$ operations in the group $`G`$, where $`|a_i|`$ denotes the length of the word $`a_i`$ in the chosen generators of the group $`G_i`$.
### 1.2 A new protocol.
In this subsection we define a new group-theoretical two party key agreement protocol that can be viewed as a non-commutative generalization of the Diffie-Hellman protocol (see ).
Let $`G`$ be a group acting on a set $`X`$ so that given $`(x,g)X\times G`$ the image $`x^g`$ of $`x`$ with respect to $`g`$ can be efficiently computed. Two parties $`A`$ and $`B`$ going to choose a secret common key from $`X`$, fix publically subgroups $`G_A,G_B`$ of the group $`G`$ and two words
$$W_A(u_A,u_B)=u_A^{a_{1,1}}u_B^{b_{1,1}}\mathrm{}u_A^{a_{1,m_1}},W_B(u_A,u_B)=u_B^{b_{2,1}}u_A^{a_{2,1}}\mathrm{}u_B^{b_{2,m_2}}$$
of the free group $`F_2`$ with two free generators $`u_A,u_B`$ such that
1. $`m_1,m_2`$, $`a_{i,j},b_{i,j}`$ for all $`i,j`$, and $`a_{1,m_1}0`$, $`b_{2,m_2}0`$,
2. $`W_A(g_A,g_B)=W_B(g_A,g_B)`$ for all $`(g_A,g_B)G_A\times G_B`$.
The protocol begins with the choice of a publically known element $`x_0X`$ and the secret elements $`g_AG_A`$ by the party $`A`$ and $`g_BG_B`$ by the party $`B`$. Then during the communications the party $`A`$ performs the following:
1. Set $`K_A=x_0`$.
2. For $`i=1,\mathrm{},m_11`$ send $`K_A^{g_{A}^{}{}_{}{}^{a_{1,i}}}`$ and receive $`K_A:=K_A^{g_{A}^{}{}_{}{}^{a_{1,i}}g_{B}^{}{}_{}{}^{b_{1,i}}}`$.
3. Set $`K_A:=K_A^{g_{A}^{}{}_{}{}^{a_{1,m_1}}}`$.
The communications of the party $`B`$ are defined similarly. Thus at the end of the communication process due to condition (W2) the parties $`A`$ and $`B`$ have the common key
$$K_A=x_0^{W_A(g_A,g_B)}=x_0^{W_B(g_A,g_B)}=K_B.$$
For $`X=_p^{}`$ with $`p`$ being a prime, $`G=G_A=G_B`$ being the group $`_{p1}^{}Aut(_p^{})`$ and $`W_A(u_A,u_B)=u_Bu_A`$, $`W_B(u_A,u_B)=u_Au_B`$ we come to the Diffie-Hellman protocol.
This scheme can be easily realized for a solvable group $`G`$ with bounded length $`n`$ of the derived series of $`G`$. For example, one can take $`G_A=G_B=G`$ and choose the words $`W_A=W_{A,n}`$ and $`W_B=W_{B,n}`$ by induction on $`n`$ as follows. If $`n=1`$, then the group $`G`$ is abelian and so conditions (W1) and (W2) are satisfied for
$$W_{A,1}(u_A,u_B)=u_Bu_A,W_{B,1}(u_A,u_B)=u_Au_B.$$
For $`n2`$ the commutator $`[g,h]=g^1h^1gh`$ with arbitrary $`g,hG`$ belongs to the derived subgroup $`G^{}=[G,G]`$ of $`G`$ (the derived length of $`G^{}`$ equals $`n1`$). Assume by induction that conditions (W1) and (W2) are satisfied for the words $`W_{A,n1}`$ and $`W_{B,n1}`$. Then a straightforward checking shows that these conditions are also satisfied, for example, for the words
$$W_{A,n}=W_{A,n1}([u_B,u_A],[u_A^1,u_B^1]),W_{B,n}=W_{B,n1}([u_B,u_A],[u_A^1,u_B^1]).$$
This follows from the fact that the length (the number of letters) of the word $`W_{A,n}`$ (as well as $`W_{B,n}`$) equals $`24^{n1}`$ which one can verify by induction on $`n1`$. More generally, to define $`W_{A,n}`$ and $`W_{B,n}`$ one can choose arbitrary words $`W_1,W_2,W_3,W_4W_X`$ where $`X=\{u_A,u_B\}`$ and $`W_X`$ is the set of all words in the alphabet $`X^\pm `$, and use $`[W_1,W_2]`$ and $`[W_3,W_4]`$ instead of $`[u_A,u_B]`$ and $`[u_B^1,u_A^1]`$ respectively. Certainly, to provide condition (1) one should guarantee that the words $`W_{A,n1}(u_A,u_B)`$ (resp. $`W_{B,n1}(u_A,u_B)`$) and $`W_2`$ (res. $`W_4`$) must be terminated to $`u_A`$ (resp. $`u_B`$). To avoid triviality we also should take $`W_1,\mathrm{},W_4`$ so that $`W_{A,n}`$ and $`W_{B,n}`$ would be nonidentity elements in the underlying free group.
Clearly, any realization of the above protocol is based on identities of the group $`G`$. In addition to commutator identities for solvable groups (see above) one can also use the identity $`x^m=1`$ (that holds in any finite group the order of which is a divisor of $`m`$, and in the Burnside groups). In this case we can choose as the words $`W_A`$ and $`W_B`$ the prefix and the inverse of the suffix of the word $`(u_Au_B)^m`$, respectively, so that the prefix is terminated to $`u_A`$. In fact, as it was proved by B.Neumann any variety of groups can be given by a collection of identities such that the first of them is of the form $`x^m=1`$ with $`m`$ being a nonnegative integer, whereas the other ones are the elements of the commutant of the underlying free group (see ).
We complete the subsection by making two remarks on the above protocol. First, the set $`X`$ must be of superpolynomial size, for otherwise the key agreement scheme can be broken in polynomial time by the known permutation group theory technique (see ). Second, the words $`W_A`$ and $`W_B`$ must be chosen so that the number of elements $`W_A(g_A,g_B)=W_B(g_A,g_B)`$ with $`g_A,g_BG`$ would contain at least two elements.
### 1.3 On the security of the protocols.
In the above protocols we assume that all groups are given explicitly, e.g. by sets of generators, so that the group operations can be performed efficiently. Then the security of the first protocol is based on the intractability of the following problem (see ).
Subgroup Conjugation Search Problem (SCSP). Given a group $`G`$, subgroups $`H_1,H_2`$ of $`G`$, and two elements $`f,gH_1`$, find an element $`hH_2`$ such that $`f=h^1gh`$, provided that at least one such $`h`$ exists.
As usually in the cryptography, an efficient algorithm solving SCSP would break the protocol (but to break the protocol it is not necessary to solve SCSP). Such an algorithm does exist for $`G=GL(n,𝔽_q)`$ where $`n`$ is a natural number, $`𝔽_q`$ is a finite field of the order $`q`$, and the subalgebra $`A(H_2)`$ of the full matrix algebra $`Mat_n(𝔽_q)`$ generated by the group $`H_2`$ is such that
$$A(H_2)G=H_2.$$
Then for arbitrary $`H_1`$ the problem SCSP can be solved in probabilistic polynomial time (in $`n`$ and in $`\mathrm{log}q`$) by the linear algebra technique, provided that $`n`$ is less than $`q/2`$. Indeed, in this case the solution of the linear system $`hfgh=0`$ with respect to $`hA(H_2)`$ is an element of $`H_2`$ with a great probability. (From it follows that in this case the problem SCSP can be solved efficiently even by a deterministic algorithm.)
It seems that the problem SCSP remains difficult when $`G`$ is restricted to subgroups of the group $`GL(V,R)`$ of all invertible $`R`$-linear transformations of the free $`R`$-module $`V`$ where $`R`$ is a finite commutative ring. To see this we consider the Linear Transporter Problem on the intractability of which the second protocol is based.
Linear Transporter Problem (LTP). Let $`R`$ be a commutative ring, $`V`$ be an $`R`$-module and $`GGL(V,R)`$. Given $`uV`$ and $`vu^G=\{u^g:gG\}`$ find $`gG`$ such that $`v=u^g`$.
A special case of (LTP) is the Discrete Logarithm Problem. Indeed, take $`V=_p^{}`$ with $`p`$ being a prime. Then $`V`$ can be considered as an one-dimensional module over the ring $`R=End(V)_{p1}`$ (with respect to taking the power $`vv^n`$ where $`vV`$, $`n_{p1}`$). Choosing $`u`$ to be a generator of the group $`V`$ we come to the Discrete Logarithm Problem.
Preserving the notation of LTP set $`T(V)=\{T_v:xx+v,v,xV\}`$ to be the translation group of the $`R`$-module $`V`$. Then obviously
$$v=u^gT_v=g^1T_ug,u,vV,gGL(V,R).$$
So the problem LTP is the special case of the problem SCSP with $`G=AGL(V,R)`$, $`H_1=T(V)`$ and $`H_2=GL(V,R)`$. (Here $`AGL(V,R)=T(V)GL(V,R)`$ is the group of all affine transformations of $`V`$.) This shows that SCSP is at least as hard as LTP. In particular, this construction gives us a family of groups for which the problem SCSP turns to be at least as hard as the Discrete Logarithm Problem. A general technique to construct groups of this kind will be given in Section 3.
## 2 Homomorphic cryptosystems over groups
### 2.1 A general scheme.
A homomorphic cryptosystem is a probabilistic public key scheme (in the sense of ) in which the spaces of plaintext messages and ciphertexts are groups $`H_k`$ and $`G_k`$ respectively, depending on a security parameter $`k`$ and such that its decryption function
$$f_k:G_kH_k$$
(2)
is an epimorphism for all $`k`$. Usually, in a homomorphic cryptosystem the public key includes generator sets $`X_k`$ and $`Y_k`$ of the groups $`G_k`$ and $`H_k`$, and some set $`R_kX_k`$ such that $`Y_kf_k(R_k)=\{f_k(g):gR_k\}`$. Besides, it is assumed that there are publically known $`k^{O(1)}`$-algorithms to solve the following problems:
1. given two elements $`a,b`$ of $`G_k`$ (resp. $`H_k`$) find the element $`ab^1`$,
2. given $`yY_k`$ find an element of the set $`R_kf_k^1(y)`$,
3. generate a random element of the group $`ker(f_k)`$
where sizes of all elements are assumed to be at most $`k`$. Under these assumptions the encryption can be performed in time $`k^{O(1)}`$ as follows. First, given a message $`h=y_1\mathrm{}y_mH_k`$ with $`y_iY_k`$ and $`m`$ being a natural number at most $`k^{O(1)}`$, Bob computes in time polynomial in $`k`$ an element $`r=r_1\mathrm{}r_mG_k`$ such that $`r_iR_k`$ and $`f_k(r_i)=y_i`$ for all $`i`$. Second, Bob mixes $`r`$ with random elements $`g_1,\mathrm{},g_{m+1}G_k`$ belonging to the kernel of the homomorphism $`f_k`$ and outputs the element $`g=g_1r_1g_2\mathrm{}g_mr_mg_{m+1}`$ as the ciphertext of $`h`$. Alice being able to compute $`f_k`$ efficiently performs the decoding as follows:
$$f_k(g)=f_k(g_1r_1g_2\mathrm{}g_mr_mg_{m+1})=f_k(r_1)\mathrm{}f_k(r_m)=y_1\mathrm{}y_m=h.$$
The key point of such a system is to choose a presentation of the group $`G_k`$ and the epimorphism $`f_k`$ in order to provide the inverse of $`f_k`$ to be a trapdoor function. The exact definition of homomorphic public-key cryptosystems and a survey of constructions can be found in .
One way to implement the general concept of a homomorphic cryptosystem is to take $`G_k`$ to be a subgroup of a certain group $`F`$ such that the group operations in $`F`$ can be performed in time polynomial in the size of operands. In the cryptosystems from and the group $`F`$ was taken as a free product of abelian groups and a modular group, respectively. In these cryptosystems the restriction of the mapping $`f_k`$ to the set $`R_k`$ was known publically and one can produce efficiently random $`k^{O(1)}`$-size elements of the group $`ker(f_k)`$. In fact, the security of these cryptosystems was based on the difficulty of the membership problem (see below) for special subgroups of the group $`G_k`$. In the next subsection we present a new homomorphic public-key cryptosystem of this kind (but with another trapdoor).
### 2.2 A new homomorphic scheme.
Let $`H=Y;`$ be a finitely presented group generated by the set $`Y`$ of cardinality $`k2`$ with $`W_Y`$ as the set of relations. As the group $`F`$ mentioned above we take the free group $`Y`$. For a permutation $`\sigma Sym(Y)`$ denote by $`\phi _\sigma `$ the automorphism of the group $`F`$ induced by $`\sigma `$. Set
$$X=X_\sigma =\{\phi _\sigma ^1(r_yyr_y^{}):yY\}$$
(3)
where $`r_y`$ and $`r_y^{}`$ are randomly chosen words of size $`O(k)`$ belonging to the set $`W_{}W_Y`$. Then $`G=X`$ is a subgroup of the group $`F`$. Moreover, the mapping $`f_\sigma :GH`$ defined by a commutative diagram
$$\begin{array}{ccc}G& \stackrel{f_\sigma }{}& H\\ id_G& & \rho \\ F& \stackrel{\phi _\sigma }{}& F\end{array}$$
(4)
where $`\rho :FH`$ is the epimorphism induced by the mapping $`id_Y`$, is an epimorphism such that given $`xX`$ we have (see (3)):
$$f_\sigma (x)=\rho (\phi _\sigma (x))=\rho (\phi _\sigma (\phi _\sigma ^1(r_yyr_y^{})))=\rho (r_yyr_y^{})=\rho (r_y)\rho (y)\rho (r_y^{})=\rho (y)=y$$
where $`y`$ is the element of $`Y`$ for which $`x=\phi _\sigma ^1(r_yyr_y^{})`$ (see (3)). In particular, $`f_\sigma (X)=Y`$ and the restriction of $`f_\sigma `$ to $`X`$ is a bijection. This enables us to construct a homomorphic cryptosystem as follows.
Secret Key: the permutation $`\sigma Sym(Y)`$.
Public Key: a natural number $`k2`$, a group $`H=Y;`$ with $`|Y|=k`$, a subgroup $`G=X_\sigma `$ of the free group $`F=Y`$, and the bijection $`f:X_\sigma Y`$ coinciding with the restriction of the homomorphism $`f_\sigma `$ to $`X_\sigma `$.
Encryption: a message $`M=y_{i_1}\mathrm{}y_{i_t}H`$ where $`y_{i_j}Y^\pm `$, is encrypted by the element
$$E(M)=f^1(s_1y_{i_1}s_1^{})\mathrm{}f^1(s_ty_{i_t}s_t^{})G$$
where $`s_i`$ and $`s_i^{}`$ are random words of the set $`W_{}W_Y`$ of size $`O(k)`$, and for a word $`w=\mathrm{}y\mathrm{}W_Y`$ we set $`f^1(w)=\mathrm{}f^1(y)\mathrm{}`$.
Decryption: a ciphertext $`C=y_{i_1}\mathrm{}y_{i_t}GF`$ where $`y_{i_j}Y^\pm `$, is decrypted to $`D(C)=y_{i_1}^\sigma \mathrm{}y_{i_t}^\sigma H`$.
To prove the correctness of the decryption we note that $`f=(f_\sigma )|_{_{X_\sigma }}`$, $`f_\sigma =\rho (\phi _\sigma )|_G`$, and $`\phi _\sigma (y)=y^\sigma `$ for all $`yY`$ (see (4)). Since obviously $`\phi _\sigma ^1=\phi _{\sigma ^1}`$, we have
$$D(E(y_{i_1}\mathrm{}y_{i_t}))=D(f^1(s_1y_{i_1}s_1^{})\mathrm{}f^1(s_ty_{i_t}s_t^{}))=D(\phi _{\sigma ^1}(s_1y_{i_1}s_1^{}\mathrm{}s_ty_{i_t}s_t^{}))=$$
$$\phi _{\sigma ^1}(s_1y_{i_1}s_1^{})^\sigma \mathrm{}\phi _{\sigma ^1}(s_ty_{i_t}s_t^{})^\sigma =s_1y_{i_1}s_1^{}\mathrm{}s_ty_{i_t}s_t^{}=y_{i_1}\mathrm{}y_{i_t}.$$
Clearly, that both encryption and decryption algorithms are polynomial-time in the size of the input words.
The security of the homomorphic cryptosystem will be discussed in the next subsection. Here we only make several remarks on the possible implementations. First, we note that it is not necessary to work with words; instead of this one can use a matrix representation of the group $`F`$ (see ). Next, to choose the set $`Y`$ so that $`|Y|=k`$, one can take any set $`S`$ of generators of $`H`$ and add to it $`k|S|`$ random elements of $`H`$ whenever $`|S|<k`$. Finally, as in Section 1 any implementation of the above cryptosystem must be supported by sufficiently large class of candidates for groups $`H`$. We will return to this problem in Section 3.
### 2.3 On the security of homomorphic schemes.
Concerning the security of the homomorphic cryptosystem suppose first that the order of the group $`H`$ is at most $`k^{O(1)}`$ (e.g. such an assumption was done in ). Then using the generator set $`Y`$ of $`H`$ one can list all the elements $`h_1,\mathrm{},h_m`$ of this group in time $`k^{O(1)}`$ and then to find within the same time a set $`\{g_1,\mathrm{},g_m\}`$ of distinct representatives of right cosets of $`G`$ by $`G_\sigma =ker(f_\sigma )`$ (one can set $`g_i=f^1(h_i)`$ for all $`i`$). Now if an adversary Charlie could recognize efficiently the elements of $`G`$ belonging to $`G_\sigma `$, then he would efficiently compute $`f_\sigma (g)`$ for all $`gG`$ due to the formulae
$$f_\sigma (g)=f_\sigma (g_i)gg_i^1G_\sigma $$
where $`i\{1,\mathrm{},m\}`$. Thus in this case the security of our cryptosystem is based on the intractability of the following problem:
Membership Testing (MT). Given a group $`F`$ and its subgroup $`G`$ test whether a given $`gF`$ belongs to $`G`$.
Suppose now that the order of $`H`$ to be arbitrary. Then a quite natural way to break the cryptosystem is to find an expression of any $`gG`$ in the terms of generators belonging to the set $`X_\sigma `$ (the attack of this kind was considered in ). Indeed, if Charlie could find efficiently for any element $`gG`$ an expression $`g=x_1\mathrm{}x_m`$ where $`x_iX_\sigma ^\pm `$ for all $`i`$, then he would efficiently compute $`f_\sigma (g)`$ due to formulae
$$f_\sigma (g)=f_\sigma (x_1)\mathrm{}f_\sigma (x_m)=f(x_1)\mathrm{}f(x_m)$$
(we recall that the bijection $`f:X_\sigma Y`$ is given publically). Thus in this case we come to the presentation problem (see ). The MT problem and the presentation problem are closely related each to other (but generally could be not polynomial-time equivalent) and one can combine them in the following well-known problem of computational group theory (see ).
Constructive Membership Testing (CMT). Given a group $`F`$ and its subgroup $`G`$ generated by a set $`X`$ find an expression of a given $`gF`$ as a word in $`X`$, or determine that $`gG`$.
Last two decades a great attention was paid to CMT with different presentations of the group $`G`$. For example, if $`F`$ is a subgroup of the symmetric group of degree $`n1`$, then the CMT can be solved in time $`n^{O(1)}`$ by the sift algorithm (see e.g. ). In the case of groups $`F=GL(n,𝔽)`$ where $`𝔽`$ is an algebraic number field, there exists an effective Las Vegas algorithm solving CMT . However, for $`n=1`$ and $`𝔽`$ being a finite field, CMT is nothing else but the the Discrete Logarithm Problem. In it was conjectured that CMT is difficult whenever the group $`G`$ either involves a large abelian group as a quotient of a normal subgroup or has nonabelian composition factors which require large degree permutation representations. Finally, the problem becomes much more difficult if we take $`F=GL(n,R)`$ the group of $`n\times n`$ invertible matrices over a ring $`R`$. In this case the problem is undecidable for $`n=4`$ and $`R=`$ (see ).
## 3 Cryptographical generation of groups
### 3.1 A general scheme.
We begin with a general scheme to construct a vast family of groups and homomorphisms supporting both key agreement protocols of Section 1 and homomorphic cryptosystems of Section 2. Let $`𝒢`$ be a class of groups closed with respect to a set $`𝒪`$ of group-theoretical operations of different arities (like direct or wreath products). For an integer $`s1`$ we denote by $`𝒪_s`$ a set of all operations of arity $`s`$ belonging to $`𝒪`$. For a set $`𝒢_0𝒢`$ we define recursively a class $`𝒫(𝒢_0,𝒪)`$ of pairs $`(G,T)`$ where $`G𝒢`$ and $`T`$ is a rooted labeled tree, as follows:
Base of recursion: any pair $`(G,T)`$ with $`G𝒢_0`$ and $`T`$ being the one-point tree with root labeled by $`G`$, belongs to $`𝒫(𝒢_0,𝒪)`$.
Recursive step: given pairs $`(G_1,T_1),\mathrm{},(G_s,T_s)𝒫(𝒢_0,𝒪)`$ and an operation $`o𝒪_s`$, the class $`𝒫(𝒢_0,𝒪)`$ contains the pair $`(G,T)`$ where $`G=o(G_1,\mathrm{},G_s)`$ and $`T`$ is the tree obtained from $`T_1,\mathrm{},T_s`$ by adding a new root labeled by $`o`$ and the sons being the roots of $`T_1,\mathrm{},T_s`$.
Let $`(G,T)𝒫(𝒢_0,𝒪)`$. Then obviously $`G𝒢`$ and the derivation tree $`T`$ of $`G`$ provides the constructive proof for this membership. The group $`G`$ is uniquely determined by $`T`$ and we call it the group associated with $`T`$. The fact, that a derivation tree is an ordinary rooted tree the leaves and the internal vertices of which are labeled by elements of $`𝒢_0`$ and $`𝒪`$ respectively, enables us to choose a random derivation tree of a fixed size.
Suppose from now on that all the groups of $`𝒢`$ are given in a certain way (e.g., one can take as $`𝒢`$ a class of matrix groups given by generator sets). We assume also that for each operation $`o𝒪_s`$ and groups $`G_1,\mathrm{},G_s𝒢`$, the size $`L(G)`$ of the presentation of the group $`G=o(G_1,\mathrm{},G_s)`$ is at most $`O(L)`$ where $`L=_{i=1}^sL(G_i)`$ and the group $`G`$ can be constructed from $`G_1,\mathrm{},G_s`$ in time $`L^{O(1)}`$. Let us define a size $`L(T)`$ of a derivation tree $`T`$ to be the sum of the sizes of all labels of $`T`$; thus $`L(T)`$ includes the sizes of the groups assigned to the leaves of $`T`$ together with the number of edges of $`T`$. Then the size of any pair $`(G,T)𝒫(𝒢_0,𝒪)`$ is $`O(L(T))`$, and the knowledge of $`T`$ enables us to find $`G`$ in time polynomial in $`L(T)`$.
One of the problems arising in constructions of group-theoretical public key cryptosystems is to find an efficient algorithm to produce a random group (or a collection of groups) belonging to a special class $`𝒢`$ and with a given size $`L`$ of the presentation. Such a group $`G`$ must be equipped with a private key providing an efficient solution of a certain computational problem for $`G`$ that is supposedly difficult in the class $`𝒢`$ without knowledge of a private key. Our approach to the above problem is to choose an appropriate class $`𝒢_0`$ of groups, a set $`𝒪`$ of group-theoretical operations, and then to generate instances for the cryptosystem in question as follows:
Step 1: given a security parameter $`L`$ choose randomly groups $`G_1,\mathrm{},G_t𝒢_0`$, such that $`_{i=1}^tL(G_i)=O(L)`$;
Step 2: choose randomly a rooted labeled tree $`T`$ of size $`O(L)`$ and with $`t`$ leaves being labeled by $`G_1,\mathrm{},G_t`$;
Step 3: compute the group $`G`$ associated with $`T`$ (i.e. $`(G,T)𝒫(𝒢_0,𝒪)`$);
Step 4: output the group $`G`$ as a public key and the labeled tree $`T`$ as a secret key.
Denote by $`𝒢^{}`$ the class of groups $`G`$ such that $`(G,T)𝒫(𝒢_0,𝒪)`$ for some labeled tree $`T`$. Then the secrecy of the key $`T`$ is based on the intractability of the following problem: given $`G𝒢^{}`$ find a derivation tree $`T`$ associated with $`G`$. A special case of this problem will be considered in Section 3.3.
For a homomorphic cryptosystem the above scheme is not sufficient because together with the group $`G`$ we have to provide a group $`H`$ and a secret homomorphism $`f:GH`$. To this end suppose that each group $`G𝒢_0`$ is equipped with a set $`M(G)`$ of pairs $`(H,f)`$ where $`H𝒢_0`$ and $`f:GH`$ is a homomorphism. We also assume that given homomorphisms $`f_i:G_iH_i`$ with $`G_i,H_i𝒢^{}`$ for $`i=1,\mathrm{},s`$, and an operation $`o𝒪_s`$ there exists an efficiently computed homomorphism $`f:GH`$ where $`G=o(G_1,\mathrm{},G_s)`$ and $`H=o(H_1,\mathrm{},H_s)`$ such that $`f|_{G_i}=f_i`$ for all $`i`$ (here we suppose in addition that $`G_i`$ is a subgroup of $`G`$). This homomorphism is denoted by $`o(f_1,\mathrm{},f_s)`$. In this notation the set $`(𝒢_0,𝒪)`$ of instances $`(G,f)`$ for a homomorphic cryptosystem can be defined recursively as follows:
Base of recursion: any pair $`(G,f)`$ with $`G𝒢_0`$ and $`fM(G)`$ belongs to the set $`(𝒢_0,𝒪)`$;
Recursion step: given pairs $`(G_1,f_1),\mathrm{},(G_s,f_s)(𝒢_0,𝒪)`$ and an operation $`o𝒪_s`$, the class $`(𝒢_0,𝒪)`$ contains the pair $`(G,f)`$ where $`G=o(G_1,\mathrm{},G_s)`$ and $`f=o(f_1,\mathrm{},f_s)`$.
We observe, that in the process of constructing the homomorphism $`f:GH`$ we also produce the derivation trees of the groups $`G`$ and $`H`$. A realization of these general schemes in finite matrix groups will be considered in the next subsection.
### 3.2 Generating matrix groups.
Let us define the classes $`𝒢_0,𝒢`$ of groups and the set $`𝒪`$ of operations. First, we set
$$𝒢=_n_R\{G:G\text{is a subgroup of}GL(n,R)\}$$
where $`n`$ and $`R`$ run over natural numbers and finite commutative rings respectively. Thus any $`G𝒢`$ is a group of $`n\times n`$ invertible matrices with entries belonging to $`R`$ for some $`n`$ and some finite commutative ring $`R`$. We recall that any such ring is a direct sum of local commutative rings and each of the latter can be described via appropriate Galois ring: the Galois ring $`GR(p^m,r)`$ of characteristic $`p^m`$ and rank $`r`$ is $`_{p^m}[x]/(f)`$ where $`f_{p^m}[x]`$ is a monic polynomial of degree $`r`$ whose image in $`_p[x]`$ is irreducible (see ).
###### Proposition 3.1
Let $`R`$ be a finite commutative local ring of characteristic $`p^m`$ and $`𝔽=GF(p^r)`$ the residue field of $`R`$. Then
1. $`R^\times =𝒯\times (1_R+Rad(R))`$ where $`𝒯`$ is a cyclic group isomorphic to $`𝔽^\times `$,
2. the subring $`R_0`$ of $`R`$ generated by $`𝒯`$ is a Galois ring $`GR(p^m,r)`$,
3. $`R`$ is a homomorphic image of the ring $`R_0[X_1,\mathrm{},X_t]`$ where $`t`$ is the minimal size of a generator set of the radical of $`R`$.
###### Proposition 3.2
Let $`p`$ be a prime and $`m,r`$ be natural numbers. Then
1. there exists the unique up to isomorphism Galois ring $`GR(p^m,r)`$ of characteristic $`p^m`$ and rank $`r`$,
2. each element $`rGR(p^m,r)`$ is uniquely represented in the form $`r=_{i=0}^{m1}t_ip^i`$ where $`t_i𝒯\{0\}`$ for all $`i`$,
3. given $`\overline{\sigma }Aut(𝔽)`$ the mapping $`r_{i=0}^{m1}t_i^\sigma p^i`$ where $`\sigma `$ is the automorphism of the group $`𝒯`$ induced by $`\overline{\sigma }`$ (see statement (1) of Proposition 3.1), is an automorphism of $`GR(p^m,r)`$.
Due to statements (2),(3) of Proposition 3.1 and statements (2) of Proposition 3.2 a representation of the finite commutative ring $`R`$ (resp., the group $`G`$) can be chosen to be polynomial in $`\mathrm{log}(|R|)`$ (resp. in $`n`$ and $`\mathrm{log}(|R|)`$). We also admit a hidden representation of $`R`$ in which the decomposition in local summands is not presented explicitly, for example the ring of residues modulo an integer can be completely given by indicating this integer.
We define a set $`𝒢_0𝒢`$ to be a class of classical simple (including abelian) subgroups $`G`$ of the groups $`GL(n,𝔽)`$ where $`n`$ and $`𝔽`$ is a finite field. Any such group $`G𝒢_0`$ is given by a set of generators so that the Membership Testing Problem can be solved in time polynomial in $`n`$ and in the bit size of $`𝔽`$. (Indeed, any nonabelian classical matrix group can be given together with a suitable matrix representation which can be used for testing membership; for an abelian group of a prime order $`p`$ one can use, e.g. the two-dimensional representation
$$_p^+GL(2,p),x\left(\begin{array}{cc}1& x\\ 0& 1\end{array}\right)$$
(5)
which gives a trivial membership testing algorithm). In fact, it is not necessary that $`𝒢_0`$ contains all classical groups; one can form $`𝒢_0`$ from the group of special types, e.g. $`PSL(n,𝔽)`$ or something like that. Since the elements of $`𝒢_0`$ are parametrized by the tuples of naturals, one can efficiently choose a random group $`G𝒢_0`$ with a given size $`L(G)`$ of presentation.
The choice of the set $`𝒪`$ of operations was inspired by the Aschbacher theorem on classifying maximal subgroups of classical groups. Let us describe the operations.
Changing the underlying ring. Let $`R`$ be a finite commutative ring and $`R^{}`$ be an extension of $`R`$. Then the natural monomorphism
$$\phi :GL(n,R)GL(n,R^{})$$
gives an unary operation in $`𝒢`$ taking $`G𝒢`$ to $`\phi (G)`$. This operation can be performed efficiently whenever e.g. the embedding $`R`$ to $`R^{}`$ is given explicitly and the number $`d=[(R^{})^+:R^+]`$ is small. Another example is the extension of $`_m`$ to $`_m^{}`$ where $`m`$ is a divisor of $`m^{}`$. Conversely, any embedding of the ring $`R^{}`$ into the ring $`Mat(d,R)`$ induces the natural monomorphism
$$\phi ^{}:GL(n,R^{})GL(nd,R)$$
taking a matrix of $`GL(n,R^{})`$ to the block matrix of $`GL(nd,R)`$ with $`d^2`$ blocks of size $`n`$. (Such a situation arises e.g. when $`R^{}`$ is a field of the order $`q^d`$ and $`R`$ is its subfield of the order $`q`$, or when $`R^{}`$ is isomorphic to the direct sum of $`d`$ copies of $`R`$.) This produces another unary operation in $`𝒢`$ taking $`G𝒢`$ to $`\phi ^{}(G)`$. In order not to blow up the representation one should assume that $`d`$ is small. In both cases the isomorphism type of the group $`G`$ (as an abstract group) does not change, but the operations change it as a linear group. In fact, our constructions start with matrix groups over a finite field $`𝔽`$. To pass to rings one can use standard extensions with $`R=𝔽`$ and $`R^{}=Mat(m,R)`$, and also with $`R=Mat(n,p)`$ and $`R^{}=Mat(m,_{p^d})`$ with a prime $`p`$.
Direct products. Suppose that groups $`G_1,\mathrm{},G_s𝒢`$ are such that $`G_iGL(n_i,R)`$ where $`n_i`$ and $`R`$ is a finite commutative ring. Then
$$G=G_1\mathrm{}G_sGL(n,R)$$
where $`n=_{i=1}^sn_i`$, and we obtain an $`s`$-ary operation in $`𝒢`$. A set of generators for the group $`G`$ can be efficiently constructed from the generating sets for $`G_1,\mathrm{},G_s`$ by means of the Kronecker product of the corresponding matrices. When $`R`$ is a field the group $`G`$ is irreducible iff so are the groups $`G_1,\mathrm{},G_s`$. (A matrix group $`G`$ is called irreducible if the underlying linear space contains no nontrivial $`G`$-invariant subspaces.)
Similarly, if $`m=n_i`$, $`G_iG_i^{}=\{I_m\}`$ and $`G_i^{}`$ normalizes $`G_i`$ where $`i=1,\mathrm{},s`$ and $`G_i^{}`$ is the group generated by $`G_j`$, $`ji`$, then $`G_1\times \mathrm{}\times G_s`$ is a subgroup of $`GL(m,R)`$ which gives one more $`m`$-ary operation.
Wreath products. The wreath product $`G\mathrm{\Gamma }`$ of a group $`G`$ and a permutation group $`\mathrm{\Gamma }Sym(m)`$ is defined to be the semidirect product of the $`m`$-fold direct product $`G^m=G\times \mathrm{}\times G`$ by the group $`\mathrm{\Gamma }`$ acting on $`G^m`$ via coordinatewise permutations. If $`GGL(n,R)`$, then the group $`G\mathrm{\Gamma }`$ has two natural linear representations obtained from the natural monomorphisms
$$G^mGL(nm,R),G^mGL(n^m,R),$$
the first of which is induced by the $`m`$-fold direct sum of the underlying $`R`$-module, whereas the second one is induced by the $`m`$-fold tensor product of it. The images of the group $`G\mathrm{\Gamma }`$ are called the imprimitive and the product actions of the wreath product, respectively. Thus we obtain two more efficiently computable $`m`$-ary operations in $`𝒢`$. In the case of $`R`$ being a field the resulting groups are always irreducible whenever $`G`$ is irreducible and $`\mathrm{\Gamma }`$ is transitive. For our purpose it is enough to set $`\mathrm{\Gamma }`$ to be the symmetric group. More elaborated way could be based on the fact that any transitive group is obtained from the action of a group on the set of right cosets by some subgroup by means of right multiplications.
Conjugations. An obvious unary operation in $`𝒢`$ consists in the conjugation of a group $`GGL(n,R)`$ by means of a randomly chosen matrix from $`GL(n,R)`$. Such an operation enables us to hide the form of a generator set of the group $`G`$.
Let $`𝒪`$ be the set of the above operations and $`𝒢^{}𝒢`$ be the set of all groups $`G`$ such that $`(G,T)𝒫(𝒢_0,𝒪)`$ for some rooted labeled tree $`T`$ (see Subsection 3.1). In the following statement we consider the specializations of the problems MT (see Subsection 2.3) and LTP (see Subsection 1.3) for the class $`𝒢^{}`$. In both cases we suppose that the group $`G𝒢^{}`$ is given by a set of generators. If $`GGL(n,R)`$ for a certain $`n`$ and for a finite commutative ring $`R`$, then in the case of LTP we set $`V`$ to be the standard free $`R`$-module of dimension $`n`$ on which the group $`GL(n,R)`$ acts, whereas for MT problem we set $`F=GL(n,R)`$.
###### Lemma 3.3
Let $`G𝒢^{}`$. Then given a derivation tree of $`G`$ the problems MT and LTP can be solved in time polynomial in $`L(G)`$.
Proof. Let $`T`$ be a derivation tree of $`G`$. Then the labels of the leaves of $`T`$ are the groups $`G_1,\mathrm{},G_t𝒢_0`$. Due to the choice of $`𝒢_0`$ the problems MT and LTP can be solved for the group $`G_i`$ in time polynomial in $`L(G_i)`$ for $`i=1,\mathrm{},t`$. Since $`L(G)=L(T)^{O(1)}`$, it suffices to verify that by means of the tree $`T`$ the problems can be reduced in time $`L(T)^{O(1)}`$ to the corresponding problems for $`G_1,\mathrm{},G_t`$. For this purpose let us consider, for instance, the reduction in the case of the primitive wreath product $`G=H\mathrm{\Gamma }`$ with $`HGL(n,R)`$ and $`\mathrm{\Gamma }=Sym(m)`$ (other operations from $`𝒪`$ on groups are treated in a similar way). Then $`GGL(n^m,R)`$ and since $`T`$ is given, we know the decomposition
$$V=U\mathrm{}U(m\text{times})$$
where $`V`$ and $`U`$ are the standard $`R`$-modules for groups $`GL(n^m,R)`$ and $`GL(n,R)`$ respectively. Any element $`gG`$ can be represented as the pair $`(h,k)H^m\times Sym(m)`$ such that
$$(u_1,\mathrm{},u_m)^g=(u_{i_1}^{h_{i_1}},\mathrm{},u_{i_m}^{h_{i_m}})$$
(6)
where $`h=(h_1,\mathrm{},h_m)`$ and $`i_j=j^{k^1}`$ for $`j=1,\mathrm{},m`$. Now the permutation $`k`$ can be efficiently computed from the elements of the form $`(0_R,\mathrm{},1_R,\mathrm{},0_R)^g`$ (with $`1_R`$ being the unique nonzero component in a certain place). So the element $`h=gg_k^1`$ also can be found efficiently where $`g_k`$ is the element of $`GL(V)=GL(n^m,R)`$ corresponding to $`k`$ (this element acts on $`V`$ exactly by permuting coordinates according to $`k`$). In particular, this provides a polynomial time reduction of the MT problem for $`G`$ to the corresponding problem for $`H`$.
Next, proceeding to the LTP problem let $`vu^G`$ for some $`u,vV`$. Denote by $`D`$ the bipartite graph with parts being the multisets $`\{u_1,\mathrm{},u_m\}`$ and $`\{v_1,\mathrm{},v_m\}`$ and the edges being the pairs $`(u_i,v_j)`$ for which $`v_i(u_j)^H`$. Then from (6) it follows that there is a one to one correspondence between the matchings $`\{(u_i,v_{j_i}):i=1,\mathrm{},m\}`$ of the graph $`D`$ and the set $`\{k\mathrm{\Gamma }:v=u^g`$ with $`g=(h,k)G`$ for some $`hH^m\}`$. Since the problem of finding a matching of a bipartite graph can be solved efficiently, we see that the LTP problem for $`G`$ is polynomial time reducible to the corresponding problem for $`H`$.
A natural way to apply our construction to the key agreement protocol is to choose a random group $`G𝒢^{}`$ of a prescribed size and then choose random subgroups $`G_A`$ and $`G_B`$ of $`G`$ (see (1)). These groups can be specified by sets of generators constructed as follows:
Step 1. Let $`S`$ be the set of leaves of the derivation tree of the group $`G`$. For each $`sS`$ take random subsets $`X_A(s)`$ and $`X_B(s)`$ of the group $`H_s`$ associated with $`s`$.
Step 2. Using the natural embedding $`hg_h`$ of $`H_s`$ into $`G`$ output $`X_A=\{g_x:xX_A(s),sS\}`$ and $`X_B=\{g_x:xX_B(s),sS\}`$ as the generator sets of $`G_A`$ and $`G_B`$ respectively.
Thus, the constructing of the groups $`G_A`$ and $`G_B`$ is performed simultaneously with the constructing the group $`G`$. (In fact, all we need is the embedding of each group assigned to a leaf of the derivation tree of the group $`G`$ into $`G`$.) In this way it is possible to control some properties of the groups, for instance, to avoid the situation when $`G_A`$ centralizes $`G_B`$ (then the common key coincides with $`1_G`$ and so is not secure).
Applying our construction to design homomorphic cryptosystems is more delicate. First of all we define the set $`M(G)`$ for each group $`GGL(n,R)`$ for some $`n`$ and some finite commutative ring $`R`$ (note that this covers the case $`G𝒢_0`$ and also allows one to produce homomorphisms in one more way: replacing $`𝒢_0`$ by a bigger subclass of $`𝒢`$). Namely, any automorphism $`\sigma Aut(R)`$ induces a homomorphism
$$f_\sigma :GG^\sigma ,AA^\sigma $$
where the matrix $`A^\sigma GL(n,R)`$ is obtained from the matrix $`AGL(n,R)`$ by entry-wise applying of $`\sigma `$. To choose $`\sigma `$ we observe that $`R=_{iI}R_i`$ where each $`R_i`$ is a finite local commutative ring. Any automorphism of the residue field of the ring $`R_i`$ can be lifted to the automorphism of this ring (statement (3) of Proposition 3.2). In the representation of the Galois ring as a quotient ring of a ring of polynomials this lifting can be done efficiently. Taking any collection $`\{\sigma _i\}_{iI}`$ one can construct the automorphism $`\sigma Aut(R)`$ such that $`\sigma |_{R_i}=\sigma _i`$ for all $`i`$. The set of such automorphisms we denote by $`Aut_0(R)`$ (in the case of $`R`$ being a field this group coincides with $`Aut(R)`$). Set
$$(G)=f_0\{f_\sigma :\sigma Aut_0(R)\}$$
(7)
where $`f_0`$ is a trivial homomorphism taking any element of $`G`$ to the identity matrix of $`GL(n,R)`$. Then assuming that the ring $`R`$ is given explicitly, one can choose a random element of $`(G)`$ in time polynomial in $`L(G)`$.
To provide the recursive step in constructing homomorphisms take $`o𝒪_s`$, $`s1`$. Suppose first that $`s=1`$. Then $`o`$ is an unary operation, i.e. either it changes the underlying ring $`R`$ of a group $`GGL(n,R)`$, or $`o`$ is a conjugation. Given a homomorphism $`f:GH`$ with $`HGL(n,R)`$ we set $`o(f)`$ to be the composition $`of`$. Let now $`s>1`$ and $`f_i:G_iH_i`$ be a homomorphism with $`G_i,H_i𝒢^{}`$ for $`i=1,\mathrm{},s`$. Then there exists the natural canonical homomorphism
$$o(f_1,\mathrm{},f_s):o(G_1,\mathrm{},G_s)o(H_1,\mathrm{},H_s)$$
coinciding with $`f_i`$ on the group $`G_i`$ which in this case is a subgroup of the group $`o(G_1,\mathrm{},G_s)`$. In any case, the resulting homomorphism is efficiently computable (we recall that we represent a homomorphism by listing explicitly the images of the generators). The above discussion shows that the following statement holds.
###### Lemma 3.4
Let $`f:GH`$ be a homomorphism constructed in the above way where $`G,H𝒢^{}`$. Then given the derivation tree of $`G`$ one can find $`f(g)`$ for $`gG`$ in time polynomial in $`L(G)`$ and the size of $`g`$.
### 3.3 Secure generation.
Let us fix the classes $`𝒢_0,𝒢,𝒢^{}`$, the set $`𝒪`$ of operations and the sets $`(G)`$ for $`G𝒢_0`$ as in Subsection 3.2. Then due to Lemmas 3.3 and 3.4 one can construct groups $`G𝒢^{}`$ to realize both key agreement protocols and homomorphic cryptosystems in which the group $`G`$ and the derivation tree $`T`$ of it play the roles of public and secret keys, respectively. The security of such systems is based on the difficulty of the following problem.
Decomposition Problem. Given a group $`G𝒢^{}`$ find a derivation tree $`T`$ of $`G`$.
This problem arises in connection with a computational version of the above mentioned Aschbacher’s theorem. A number of practical algorithms (without complexity bounds) for Decomposition Problem are known (see ) but in general this problem seems to be difficult. Indeed, suppose that $`R=_m`$ where $`m=pq`$ with $`p`$ and $`q`$ being two different primes. Denote by $`G_p`$ the cyclic matrix group of the order $`p`$ in $`GL(2,p)`$ (see (5)). Similarly, the group $`G_q`$ is defined. Then $`G_p,G_q𝒢_0`$ and
$$G=G_p^{}\times G_q^{}GL(2,R)$$
where $`G_p^{}`$ and $`G_q^{}`$ are the images of the groups $`G_p`$ and $`G_q`$ with respect to the natural embeddings $`GL(2,p)`$ and $`GL(2,q)`$ into $`GL(2,R)`$. Thus the group $`G`$ can be constructed in two steps: construct the groups $`G_p^{}`$ and $`G_q^{}`$ (the operation of changing the underlying ring), and set $`G=G_1^{}\times G_2^{}`$ (the operation of the direct product). This implies that $`G𝒢^{}`$. This shows that the integer factoring problem is a special case of the Decomposition Problem.
Another strategy of Charlie could be to avoid solving the Decomposition Problem and to try solve the problems like LTP, SCSP or CMT directly. To prevent such an attack one can choose the leaves of the derivation tree of the group $`G`$ to be the groups of the exponential size with respect to $`L(G)`$. Then from the construction it follows that these groups will arise as the composition factors of $`G`$. However, for the groups with large composition factors all the problems like LTP, SCSP or CMT seem to be difficult (see Subsections 1.3 and 2.3).
We mention one more attack of Charlie for the case of a homomorphic cryptosystem. Suppose we construct in the above way the homomorphism $`f:GH`$ with $`G,H𝒢^{}`$. We call the homomorphism linear if it induces the ring homomorphism $`f^{}:A(G)A(H)`$ where $`A(G)`$ (resp. $`A(H)`$) is the subring of the underlying full matrix ring generated by $`G`$ (resp. $`H`$). For a linear homomorphism the corresponding homomorphic cryptosystem can be easily broken whenever $`GGL(n,R)`$ where $`R=_n`$ for some $`n`$ or $`R`$ is a finite field (or, more generally, a direct sum of Galois rings). Indeed, in this case Charlie can find $`f(g)`$ for $`gG`$ as follows. Take random generators $`g_1,\mathrm{},g_s`$ of the group $`G`$ and find a decomposition $`g=_{i=1}^sc_ig_i`$ with $`c_iR`$ just involving linear algebra. Then $`f(g)=_{i=1}^sc_if(g_i)`$ due to the linearity of $`f`$. To prevent this attack one can take some initial homomorphisms at the leaves of the derivation tree to be elements of the group $`Aut_0(R)`$ (see (7)). Then the constructed homomorphism is not linear in general (e.g. if $`gGL(n,𝔽)`$ with $`𝔽`$ being a field, and $`\sigma Aut(𝔽)`$, then generally $`(ag)^\sigma ag^\sigma `$).
We complete the subsection by the following statement summarizing the above discussion.
###### Theorem 3.5
Assuming that the problems LTP, SCSP, CMT for matrix groups over finite commutative rings, as well as the Decomposition Problem are intractable, a secure two-party key agreement protocol and homomorphic cryptosystem can be implemented for these groups.
One of the consequences of this theorem is that by means of it one can construct encrypted simulation of a boolean circuit of the logarithmic depth (the details can be found in ).
## Final remarks
One of the main problems in constructing homomorphic public-key cryptosystems consists in finding appropriate trapdoor functions. However, in the natural presentations of homomorphisms of algebraic structures the problem of breaking such a system is reduced to some variants of the CMT problem. On the other hand, there is the following result for matrix groups over finite fields.
###### Theorem 3.6
\[12, Theorem 6.1\] Given $`K=XGL(d,p^e)`$ where $`XGL(d,p^e)`$, there is a Las Vegas algorithm that given any $`gGL(d,p^e)`$, decides whether $`gK`$, and if $`gK`$, then the algorithm produces a straight-line program with the input $`X`$, yielding g. The algorithm uses an oracle to compute discrete logarithms in fields of characteristic $`p`$ with sizes up to $`p^{ed}`$. In case when all of those composition factors of Lie type in characteristic $`p`$ are constructively recognizable with a Discrete Log oracle <sup>1</sup><sup>1</sup>1The current list of groups of Lie type recognizable with a Discrete Log oracle is given in ; this list includes the groups of series A, B, C, D., the running time is a polynomial in the input length $`|X|d^2e\mathrm{log}p`$, plus the time required for polynomially many calls to the Discrete Log oracle.
This theorem shows that having an oracle for the Discrete Logarithm, the membership problem can be solved in probabilistic polynomial time for matrix groups over finite fields. This means that at least for homomorphic public-key cryptosystems over such groups there is a little hope to find a trapdoor function different from functions the difficulty of inversion of which is based on the intractability of the Discrete Logarithm. However, only a little is known on the computational complexity of the membership problem for matrix groups over rings. So constructions over such groups seems to be more perspective from the point of view of algebraic (non-commutative) cryptography. |
warning/0506/astro-ph0506308.html | ar5iv | text | # Fast directional continuous spherical wavelet transform algorithms
## I Introduction
Wavelet analysis has proven useful in many applications due to the ability of wavelets to resolve localised signal content in both scale and space. Many of these applications, however, are restricted to data defined in Euclidean space: the 1-dimensional line, the 2-dimensional plane and, occasionally, higher dimensions. Nevertheless, data are often measured or defined on other manifolds, such as the 2-sphere. Examples where data are measured on the sphere are found in astrophysics (e.g. ), planetary science (e.g. ), geophysics (e.g. ), computer vision (e.g. ) and quantum chemistry (e.g. ). To realise the potential benefits that may be provided by wavelets in such settings, ordinary Euclidean wavelet analysis must be extended to spherical geometry.
A number of attempts have been made to extend wavelets to the unit sphere. Discrete second generation wavelets on the sphere that are based on a multiresolution analysis have been developed . Other authors have focused on the continuous wavelet transform on the sphere. A number of works construct a solution using a harmonic approach , however these solutions suffer from the poor localisation of the spherical harmonic functions. Others adopt a tangent bundle viewpoint , thereby avoiding the necessity to define a dilation operator on the sphere. A satisfactory extension of the continuous wavelet transform to the sphere is defined by , however this construction requires an abstract dilation parameter that must satisfy a number of ad hoc assumptions. More recently, a consistent and satisfactory framework for wavelets defined on the unit sphere has been constructed and developed by . Moreover, this construction is derived entirely from group theoretic principles and inherently satisfies a number of natural requirements. We consider the continuous spherical wavelet transform (CSWT) developed in these last works. For a more detailed review of the attempts made at constructing a wavelet transform on the unit sphere see .
Current and future data-sets defined on the sphere are of considerable size. The current Wilkinson Microwave Anisotropy Probe (WMAP) data of the cosmic microwave background (CMB) contain approximately $`3\times 10^6`$ pixels on the sphere, whereas the forthcoming Planck CMB mission will generate maps with approximately $`50\times 10^6`$ pixels. Fast algorithms are therefore required to perform the CSWT on practical data-sets. A semi-fast algorithm to implement the CSWT is presented in and implemented in the YAWTb<sup>1</sup><sup>1</sup>1http://www.fyma.ucl.ac.be/projects/yawtb/ (Yet-Another-Wavelet-Toolbox) Matlab wavelet toolbox (which also makes use of the SpharmonicKit<sup>2</sup><sup>2</sup>2http://www.cs.dartmouth.edu/~geelong/sphere/). However, this implementation is restricted to an equi-angular tessellation of the sphere. The beauty of this tessellation is its simplicity and ability to be easily represented in matrix form. However, the pixels of an equi-angular tessellation are densely spaced about the poles and do not have equal areas. Other tessellations of the sphere also exist, such as those constructed to minimise some energy measure or those constructed for more practical or numerical purposes (for example the IGLOO<sup>3</sup><sup>3</sup>3http://www.mrao.cam.ac.uk/projects/cpac/igloo/ , HEALPix<sup>4</sup><sup>4</sup>4http://healpix.jpl.nasa.gov/ and GLESP<sup>5</sup><sup>5</sup>5http://www.glesp.nbi.dk/ schemes). There is thus a need for a fast implementation of the CSWT that is not tied to any particular tessellation of the sphere.
We fill this void by presenting a fast algorithm for performing the directional CSWT. The CSWT at a particular scale is essentially a spherical convolution, hence we may apply the fast spherical convolution algorithm developed by to evaluate the wavelet transform. The algorithm is posed in harmonic space and thus is independent of the underlying tessellation of the sphere, (although an iso-latitude tessellation does enable faster spherical harmonic transforms, thereby increasing the speed of the algorithm). The framework supports both non-azimuthally symmetric spherical wavelets<sup>6</sup><sup>6</sup>6Azimuthally symmetric spherical wavelets are also often referred to as axisymmetric wavelets. and a decomposition that employs anisotropic dilations, however no synthesis is possible when anisotropic dilations are incorporated. For an illustration of a spherical wavelet analysis on a practical problem of considerable size we refer the reader to our recent works to test the CMB for deviations from Gaussianity and to detect the integrated Sachs-Wolfe (ISW) effect . Both of these works involve performing 1000 Monte Carlo simulations and would not have been feasible without a fast directional CSWT algorithm.
The remainder of this paper is structured as follows. In section II we describe the CSWT in the framework presented by . For a more complete treatment of the spherical wavelet transform in this framework and the correspondence between spherical and Euclidean wavelets we recommend that the reader refer to . We also present an extension to anisotropic dilations, however in this case the basis functions are not strictly wavelets hence perfect reconstruction is not possible. Various algorithms to perform the CSWT are described and then compared in section III. An application of our implementation is demonstrated in section IV. Concluding remarks are made in section V.
## II The continuous spherical wavelet transform
To extend Euclidean wavelet analysis to spherical geometry a number of requirements must be satisfied: (i) the signals and wavelets must live fully on the unit sphere; (ii) the transform must involve local dilations of some kind on the unit sphere; and (iii) the spherical wavelet transform should reduce locally to the Euclidean transform on the tangent plane (i.e. the Euclidean limit must be satisfied) . The final requirement is intuitively obvious; the sphere is asymptotically flat, hence any spherical wavelet transform should match the planar Euclidean transform on small scales, or equivalently, for a large radius of curvature.
The spherical wavelet transform developed in satisfies all of these requirements, moreover each requirement naturally follows from the construction. The construction of this transform is derived entirely from group theoretical principles. However, in a recent work by this formalism is reintroduced independently of the original group theoretic formalism, in an equivalent, practical and self-consistent approach. We adopt this approach herein.
The correspondence principle between spherical and Euclidean wavelets is developed by , relating the concepts of planar Euclidean wavelets to spherical wavelets through a stereographic projection. We use the stereographic projection to define affine transformations on the unit sphere that facilitate the construction of a wavelet basis on the unit sphere. The spherical wavelet transform may then be defined as the projection on to this basis, where the spherical wavelets must satisfy the appropriate admissibility criterion to ensure perfect reconstruction.
### II-A Stereographic projection
In order to construct a correspondence between wavelets on the plane ($`^2`$) and sphere ($`S^2`$) a projection operator between the two spaces must be chosen. It is shown in that the stereographic projection is the unique unitary, radial and conformal diffeomorphism between the sphere and the plane.
The stereographic projection is defined by projecting a point on the unit sphere to a point on the tangent plane at the north pole, by casting a ray though the point and the south pole. The point on the unit sphere is mapped on to the intersection of this ray and the tangent plane (see Fig. 1). Formally, we may define the stereographic projection operator as $`\mathrm{\Pi }:\omega 𝐱=\mathrm{\Pi }\omega =(r(\theta ),\varphi )`$ where $`r=2\mathrm{tan}(\theta /2)`$, $`\omega (\theta ,\varphi )S^2`$ denotes spherical coordinates with colatitude $`\theta `$ and longitude $`\varphi `$ and $`𝐱^2`$ is a point in the plane, denoted here by the polar coordinates $`(r,\varphi )`$. The inverse operator is $`\mathrm{\Pi }^1:𝐱\omega =\mathrm{\Pi }^1𝐱=(\theta (r),\varphi )`$, where $`\theta (r)=2\mathrm{tan}^1(r/2)`$.
Following the formulation of again, we define the action of the stereographic projection operator on functions on the plane and sphere. We consider the space of square integrable functions in $`L^2(^2,\mathrm{d}^2𝐱)`$ on the plane and $`L^2(S^2,\mathrm{d}\mathrm{\Omega })`$ on the unit sphere, where $`\mathrm{d}\mathrm{\Omega }=\mathrm{sin}\theta \mathrm{d}\theta \mathrm{d}\varphi `$ is the usual rotation invariant measure on the sphere. The action of the stereographic projection operator $`\mathrm{\Pi }:sL^2(S^2,\mathrm{d}\mathrm{\Omega })p=\mathrm{\Pi }sL^2(^2,\mathrm{d}^2𝐱)`$ on functions is defined as
$$p(r,\varphi )=(\mathrm{\Pi }s)(r,\varphi )=(1+r^2/4)^1s(\theta (r),\varphi ).$$
(1)
The inverse stereographic projection operator $`\mathrm{\Pi }^1:pL^2(^2,\mathrm{d}^2𝐱)s=\mathrm{\Pi }^1pL^2(S^2,\mathrm{d}\mathrm{\Omega })`$ on functions is then
$$s(\theta ,\varphi )=(\mathrm{\Pi }^1p)(\theta ,\varphi )=[1+\mathrm{tan}^2(\theta /2)]p(r(\theta ),\varphi ).$$
(2)
The pre-factors introduced ensure that the $`L^2`$-norm of functions through the forward and inverse projections are conserved. In the Euclidean limit, the stereographic projection and inverse naturally reduce to the identity operator .
### II-B Affine transformations on the sphere
A wavelet basis is constructed on the unit sphere in section II-C by applying the spherical extension of Euclidean translations and dilations to mother wavelets defined on the unit sphere. The extension of these affine transformations to the sphere, facilitated by the stereographic projection operator, are defined here.
The natural extension of Euclidean translations on the unit sphere are rotations. These are characterised by the elements of the rotation group $`\mathrm{SO}(3)`$, which we parameterise in terms of the three Euler angles $`\rho =(\alpha ,\beta ,\gamma )`$.<sup>7</sup><sup>7</sup>7We adopt the $`zyz`$ Euler convention corresponding to the rotation of a physical body in a *fixed* co-ordinate system about the $`z`$, $`y`$ and $`z`$ axes by $`\gamma `$, $`\beta `$ and $`\alpha `$ respectively. The rotation of a square-integrable function $`sL^2(S^2,\mathrm{d}\mathrm{\Omega })`$ is defined by
$$[R(\rho )s](\omega )=s(\rho ^1\omega ),\rho \mathrm{SO}(3).$$
(3)
Dilations on the unit sphere are constructed by first lifting the sphere to the plane by the stereographic projection, followed by the usual Euclidean dilation in the plane, before re-projecting back onto the sphere. We generalise to anisotropic dilations on the sphere (a similar anisotropic dilation operator on the sphere has been independently proposed by ), however in this setting we do not achieve a wavelet basis and hence cannot synthesise our original signal. We define the anisotropic Euclidean dilation operator in $`L^2(^2,\mathrm{d}^2𝐱)`$ as
$$[d(a,b)p](x,y)=a^{1/2}b^{1/2}p(a^1x,b^1y),$$
(4)
for the non-zero positive scales $`a,b_{}^+`$. The $`a^{1/2}b^{1/2}`$ normalisation factor ensures the $`L^2`$-norm is preserved. The spherical dilation operator $`𝒟(a,b):s(\theta ,\varphi )[𝒟(a,b)s](\theta ,\varphi )`$ in $`L^2(S^2,\mathrm{d}\mathrm{\Omega })`$ is defined as the conjugation by $`\mathrm{\Pi }`$ of the Euclidean dilation $`d(a,b)`$ in $`L^2(^2,\mathrm{d}^2𝐱)`$ on the tangent plane at the north pole:
$$𝒟(a,b)=\mathrm{\Pi }^1d(a,b)\mathrm{\Pi }.$$
(5)
The norm of functions in $`L^2(S^2,\mathrm{d}\mathrm{\Omega })`$ is preserved by the spherical dilation as both the stereographic projection operator and Euclidean dilations preserve the norm of functions. Extending the isotropic spherical dilation operator defined by to anisotropic dilations, we obtain
$$[𝒟(a,b)s](\omega )=[\lambda (a,b,\theta ,\varphi )]^{1/2}s(\omega _{1/a,1/b}),$$
(6)
where $`\omega _{a,b}=(\theta _{a,b},\varphi _{a,b})`$,
$$\mathrm{tan}(\theta _{a,b}/2)=\mathrm{tan}(\theta /2)\sqrt{a^2\mathrm{cos}^2\varphi +b^2\mathrm{sin}^2\varphi }$$
and
$$\mathrm{tan}(\varphi _{a,b})=\frac{b}{a}\mathrm{tan}(\varphi ).$$
For the case where $`a=b`$ the anisotropic dilation reduces to the usual isotropic case defined by $`\mathrm{tan}(\theta _a/2)=a\mathrm{tan}(\theta /2)`$ and $`\varphi _a=\varphi `$. The $`\lambda (a,b,\theta ,\varphi )`$ cocycle term follows from the factors introduced in the stereographic projection of functions to preserve the $`L^2`$-norm. Alternatively, the cocycle may be derived explicitly to preserve the $`L^2`$-norm when the stereographic projection of functions do not have these pre-factor terms. The cocycle of an anisotropic spherical dilation is defined by
$$\lambda (a,b,\theta ,\varphi )=\frac{4a^3b^3}{\left(A_{}\mathrm{cos}\theta +A_+\right)^2},$$
(7)
where
$$A_\pm =a^2b^2\pm a^2\mathrm{sin}^2\varphi \pm b^2\mathrm{cos}^2\varphi .$$
For the case where $`a=b`$ the anisotropic cocycle reduces to the usual isotropic cocycle
$$\lambda (a,a,\theta ,\varphi )=\frac{4a^2}{[(a^21)\mathrm{cos}\theta +a^2+1]^2}.$$
Although the ability to perform anisotropic dilations is of practical use, we do not achieve a wavelet basis in this setting. In the anisotropic setting a bounded admissibility integral cannot be determined (even in the plane), thus the synthesis of a signal from its coefficients cannot be performed. This results from there being no direct means of evaluating the proper measure in the absence of a group structure. The projection of a signal onto basis functions undergoing anisotropic dilations may be performed in an analogous manner to the following discussion of the wavelet transform. However, since these basis functions are not wavelets we restrict the following discussion to isotropic dilations.
### II-C Wavelet transform
A wavelet basis on the unit sphere may now be constructed from rotations and isotropic dilations (where $`a=b`$) of a mother spherical wavelet $`\psi L^2(S^2,\mathrm{d}\mathrm{\Omega })`$. The corresponding wavelet family $`\{\psi _{a,\rho }R(\rho )𝒟(a,a)\psi ,\rho \mathrm{SO}(3),a_{}^+\}`$ provides an over-complete set of functions in $`L^2(S^2,\mathrm{d}\mathrm{\Omega })`$. The CSWT of $`sL^2(S^2,\mathrm{d}\mathrm{\Omega })`$ is given by the projection onto each wavelet basis function in the usual manner,
$$W_\psi ^s(a,\rho )_{S^2}d\mathrm{\Omega }\psi _{a,\rho }^{}(\omega )s(\omega )=\psi _{a,\rho }s,$$
(8)
where the $``$ denotes complex conjugation.
The transform is general in the sense that all orientations in the rotation group $`\mathrm{SO}(3)`$ are considered, thus directional structure is naturally incorporated. It is important to note, however, that only *local* directions make any sense on $`S^2`$. There is no global way of defining directions on the sphere<sup>8</sup><sup>8</sup>8There is no differentiable vector field of constant norm on the sphere and hence no global way of defining directions. – there will always be some singular point where the definition fails.
The synthesis of a signal on the unit sphere from its wavelet coefficients is given by
$$s(\omega )=_{\mathrm{SO}(3)}d\rho _0^{\mathrm{}}\frac{\mathrm{d}a}{a^3}W_\psi ^s(a,\rho )[R(\rho )L_\psi \psi _a](\omega ),$$
(9)
where $`\mathrm{d}\rho =\mathrm{sin}\beta \mathrm{d}\alpha \mathrm{d}\beta \mathrm{d}\gamma `$. The $`L_\psi `$ operator in $`L^2(S^2,\mathrm{d}\mathrm{\Omega })`$ is defined by the action
$$(\widehat{L_\psi f})_\mathrm{}m=\widehat{f}_\mathrm{}m/C_\psi ^{\mathrm{}}$$
(10)
on the spherical harmonic coefficients of functions $`fL^2(S^2,\mathrm{d}\mathrm{\Omega })`$, where $`C_\psi ^{\mathrm{}}`$ is defined below. The hat denotes the spherical harmonic coefficients
$$\widehat{f}_\mathrm{}m=_{S^2}d\mathrm{\Omega }Y_\mathrm{}m^{}(\omega )f(\omega )=Y_\mathrm{}mf$$
of the decomposition
$$f(\omega )=\underset{\mathrm{}=0}{\overset{\mathrm{}}{}}\underset{m=\mathrm{}}{\overset{\mathrm{}}{}}\widehat{f}_\mathrm{}mY_\mathrm{}m(\omega ).$$
We adopt the Condon-Shortley phase convention where the normalised spherical harmonics are defined by
$$Y_\mathrm{}m(\omega )=(1)^m\sqrt{\frac{2\mathrm{}+1}{4\pi }\frac{(\mathrm{}m)!}{(\mathrm{}+m)!}}P_{\mathrm{}}^m(\mathrm{cos}\theta )e^{im\varphi },$$
where $`P_{\mathrm{}}^m(x)`$ are the associated Legendre functions. Using this normalisation the orthogonality of the spherical harmonic functions is given by
$$_{S^2}d\mathrm{\Omega }Y_\mathrm{}m(\omega )Y_\mathrm{}^{}m^{}^{}(\omega )=\delta _{\mathrm{}\mathrm{}^{}}\delta _{mm^{}},$$
(11)
where $`\delta _{ij}`$ is Kronecker delta function. In order to ensure the perfect reconstruction of a signal synthesised from its wavelet coefficients, one requires the admissibility condition
$$0<C_\psi ^{\mathrm{}}\frac{8\pi ^2}{2\mathrm{}+1}\underset{m=\mathrm{}}{\overset{\mathrm{}}{}}_0^{\mathrm{}}\frac{\mathrm{d}a}{a^3}(\widehat{\psi _a})_\mathrm{}m^2<\mathrm{}$$
(12)
to hold for all $`\mathrm{}`$. A proof of the admissibility criterion is given by . Practically, it is difficult to apply (12) directly, thus a necessary (and almost sufficient) condition for admissibility is the zero-mean condition
$$C_\psi _{S^2}d\mathrm{\Omega }\frac{\psi (\omega )}{1+\mathrm{cos}\theta }=0.$$
(13)
In the Euclidean limit this condition naturally reduces to the necessary zero-mean condition for Euclidean wavelets .
### II-D Correspondence principle and spherical wavelets
The correspondence principle between spherical and Euclidean wavelets states that the inverse stereographic projection of an *admissible* wavelet on the plane yields an *admissible* wavelet on the unit sphere. This result is proved by . Hence, mother spherical wavelets may be constructed from the projection of mother Euclidean wavelets on the plane:
$$\psi (\omega )=[\mathrm{\Pi }^1\psi _^2](\omega ),$$
(14)
where $`\psi _^2L^2(^2,\mathrm{d}^2𝐱)`$ is an admissible wavelet in the plane. Directional spherical wavelets may be naturally constructed in this setting – they are simply the projection of directional Euclidean planar wavelets on to the sphere.
We give examples of three spherical wavelets: the spherical Mexican hat wavelet (SMHW); the spherical butterfly wavelet (SBW); and the spherical real Morlet wavelet (SMW). These spherical wavelets are illustrated in Fig. 2. Each spherical wavelet is constructed by the stereographic projection of the corresponding Euclidean wavelet onto the sphere, where the Euclidean planar wavelets are defined by
$$\psi _^2^{\mathrm{SMHW}}(r,\varphi )=\frac{1}{2}(2r^2)e^{r^2/2},$$
$$\psi _^2^{\mathrm{SBW}}(x,y)=xe^{(x^2+y^2)/2}$$
and
$$\psi _^2^{\mathrm{SMW}}(𝐱;𝐤)=\mathrm{Re}\left(e^{i𝐤𝐱/\sqrt{2}}e^{𝐱^2/2}\right)$$
respectively, where $`𝐤`$ is the wave vector of the SMW. The SMHW is proportional to the Laplacian of a Gaussian, whereas the SBW is proportional to the first partial derivative of a Gaussian in the $`x`$-direction. The SMW is a Gaussian modulated sinusoid, or Gabor wavelet.
A full directional wavelet analysis on the unit sphere for large data sets has previously been prohibited by the computational infeasibility of any implementation. The computational burden of computing many orientations may be reduced by using steerable wavelets, for which any continuous orientation can be computed from a small number of basis orientations . This is achieved since steerable wavelets have a limited azimuthal band limit and may thus be represented as a finite sum of trigonometric exponentials . However, in this case one must still compute the initial transform for more than one orientation, so although the computational burden is reduced, it is still significant. Moreover, we also require a fast approach for general non-steerable, directional wavelets. We address this problem in the following section by presenting a fast algorithm to perform the directional CSWT.
## III Algorithms
A range of algorithms of varying computational efficiency and numerical accuracy are presented to perform the CSWT described in section II. We implement these algorithms in Fortran 90 and subsequently compare computational complexity and typical execution time. The synthesis of a signal from its wavelet coefficients is not considered any further. Without loss of generality we consider only a single dilation (i.e. fixed $`a`$ and $`b`$).
### III-A Tessellation schemes
It is necessary to discretise both the spherical coordinates of a function defined on the unit sphere and also the Euler angle representation of the $`\mathrm{SO}(3)`$ rotation group. The fast algorithms we present are performed in harmonic space and hence are tessellation independent, provided an appropriate spherical harmonic transform is defined for the tessellation. However, the semi-fast algorithm is restricted to an equi-angular tessellation of the sphere. The various tessellation schemes adopted are defined below.
The equi-angular tessellation (also known as the equidistant cylindrical projection (ECP)) of the spherical coordinates is defined by $`𝒞=\{\theta _{n_\theta }=\frac{\pi n_\theta }{N_\theta },\varphi _{n_\theta }=\frac{2\pi n_\varphi }{N_\varphi }:0n_\theta N_\theta 1,\mathrm{\hspace{0.17em}0}n_\varphi N_\varphi 1\}`$. Let $`N_{\mathrm{pix}}=N_\theta N_\varphi `$ denote the number of pixels in the tessellation.
We also consider the HEALPix tessellation scheme since it is commonly used for astrophysical data-sets of the CMB. Pixels in the HEALPix scheme are of equal area and are located on rings of constant latitude (the latter feature enables fast spherical harmonic transforms on the pixelised grid). We refer the reader to for details of the HEALPix scheme and here just define the HEALPix grid in terms of pixel indices: $`=\{(\theta ,\varphi )_p=(\theta _p,\varphi _p):0pN_{\mathrm{pix}}1\}`$. The HEALPix resolution is parameterised by $`N_{\mathrm{side}}`$, where $`N_{\mathrm{pix}}=12N_{\mathrm{side}}^{}{}_{}{}^{2}`$. It should be noted that an exact quadrature formula does not exist for the HEALPix tessellation, thus spherical harmonic transforms are necessarily approximate. This is not the case for the ECP or other practical tessellations (e.g. IGLOO and GLESP) where exact quadrature may be performed.
The Euler angle domain of the spherical wavelet coefficients is in general arbitrary, however we use the equi-angular discretisation defined by $`_1=\{\alpha _{n_\alpha }=\frac{2\pi n_\alpha }{N_\alpha },\beta _{n_\beta }=\frac{\pi n_\beta }{N_\beta },\gamma _{n_\gamma }=\frac{2\pi n_\gamma }{N_\gamma }:0n_\alpha N_\alpha 1,\mathrm{\hspace{0.17em}0}n_\beta N_\beta 1,\mathrm{\hspace{0.17em}0}n_\gamma N_\gamma 1\}`$. Our fast algorithm, however, requires (for convenience) the tessellation $`_2=\{\alpha _{n_\alpha }=\frac{2\pi n_\alpha }{N_\alpha },\beta _{n_\beta }=\frac{2\pi n_\beta }{N_\beta },\gamma _{n_\gamma }=\frac{2\pi n_\gamma }{N_\gamma }:0n_\alpha N_\alpha 1,\mathrm{\hspace{0.17em}0}n_\beta 2N_\beta 1,\mathrm{\hspace{0.17em}0}n_\gamma N_\gamma 1\}`$, where the $`\beta `$ sampling is repeated. Evaluating $`\beta `$ over the range $`0`$ to $`2\pi `$ is redundant, covering the $`\mathrm{SO}(3)`$ manifold exactly twice. Nonetheless, the use of our fast algorithm requires this range. Such an approach is not uncommon, as also oversample a function of the sphere in the $`\theta `$ direction in order to develop fast spherical harmonic transforms on equi-angular grids. The overhead associated with our inefficient discretisation is more than offset by the fast algorithm it affords, as described in section III-E.
### III-B Direct case
The CSWT defined by (8) may be implemented directly by applying an appropriate quadrature rule. Using index subscripts to denote sampled signals, the direct CSWT implementation is given by
$`(W_\psi ^s)_{n_\alpha ,n_\beta ,n_\gamma }=`$ (15)
$`{\displaystyle \underset{p=0}{\overset{N_{\mathrm{pix}}1}{}}}\left[R({\displaystyle \frac{2\pi n_\alpha }{N_\alpha }},{\displaystyle \frac{\pi n_\beta }{N_\beta }},{\displaystyle \frac{2\pi n_\gamma }{N_\gamma }})\psi \right]_p^{}s_pw_p,`$
where the pixel sum is over all the pixels of any chosen tessellation. The weights for the ECP $`𝒞`$ grid are given by $`w_p=w_{n_\theta }=\frac{2\pi ^2\mathrm{sin}\theta _{n_\theta }}{N_\theta N_\varphi }`$, whereas the equal pixel areas of the HEALPix $``$ grid ensure the pixel weights, given by $`w_p=\frac{4\pi }{N_{\mathrm{pix}}}`$, are independent of position.
Discretisation techniques other than the plain Riemann sum used in (15) would be beneficial only if additional regularity conditions are imposed on the signal $`s`$ . It is also possible to choose other weights to achieve a better approximation of the integral. An example of a different equi-angular discretisation and a different choice for the weights is given by the sampling theorem for band-limited functions on the sphere developed by .
Evaluation of (15) requires the computation of a 2-dimensional summation, evaluated over a 3-dimensional grid. We assume the number of samples for each discretised angle, except $`\gamma `$, is of the same order $`N`$. Typically the number of samples in the $`\gamma `$ direction is of a much lower order, so we treat this term separately. The complexity of the direct algorithm is $`𝒪(N_\gamma N^4)`$.
### III-C Semi-fast case
We rederive here the semi-fast implementation of the CSWT described by and implemented using the YAWTb Matlab wavelet toolbox and SpharmonicKit. This algorithm involves performing a separation of variables so that one rotation may be performed in Fourier space. The algorithm is restricted to the equi-angular grid $`𝒞`$ (in essence pixels must only be defined on equal latitude rings, however some form of interpolation and down-sampling is then required to extract samples for equal longitudes).
Firstly, the $`\alpha `$ rotation is represented by shifting the corresponding wavelet samples to give
$`(W_\psi ^s)_{n_\alpha ,n_\beta ,n_\gamma }=`$
$`{\displaystyle \underset{n_\theta =0}{\overset{N_\theta 1}{}}}{\displaystyle \underset{n_\varphi =0}{\overset{N_\varphi 1}{}}}\left[R(0,{\displaystyle \frac{\pi n_\beta }{N_\beta }},{\displaystyle \frac{2\pi n_\gamma }{N_\gamma }})\psi \right]_{n_\theta ,n_\varphi n_\alpha }^{}s_{n_\theta ,n_\varphi }w_{n_\theta },`$
where the index $`n_\varphi `$ is extended periodically with period $`N_\varphi `$. The discrete space convolution theorem may then be applied to represent the inner summation as the inverse discrete Fourier transform (DFT) of the product of the wavelet and signal DFT samples (note that only a 1-dimensional DFT is performed in the azimuthal direction):
$`(W_\psi ^s)_{n_\alpha ,n_\beta ,n_\gamma }=`$ (16)
$`{\displaystyle \underset{n_\theta =0}{\overset{N_\theta 1}{}}}\{{\displaystyle \frac{1}{N_\varphi }}{\displaystyle \underset{k=0}{\overset{N_\varphi 1}{}}}^{}\left[R(0,{\displaystyle \frac{\pi n_\beta }{N_\beta }},{\displaystyle \frac{2\pi n_\gamma }{N_\gamma }})\psi \right]_{n_\theta ,k}`$
$`\times (s)_{n_\theta ,k}e^{\frac{i2\pi kn_\varphi }{N_\varphi }}\}w_{n_\theta },`$
where $`()_{n,k}`$ denotes 1-dimensional DFT coefficients. A fast Fourier transform (FFT) may then be applied to evaluate simultaneously all of the $`n_\alpha `$ terms of the expression enclosed in the curly braces in (16). A final summation (integral) over $`n_\theta `$ produces the spherical wavelet coefficients for a given $`n_\beta `$ and $`n_\gamma `$, for *all* $`n_\alpha `$. Applying an FFT to evaluate simultaneously one summation rapidly, reduces the complexity of the CSWT implementation to $`𝒪(N_\gamma N^3\mathrm{log}_2N)`$.
### III-D Fast azimuthally symmetric case
The fast azimuthally symmetric CSWT algorithm is posed in harmonic space, where $`\widehat{s}_\mathrm{}m=Y_\mathrm{}ms`$ are the spherical harmonic coefficients of a function $`sL^2(S^2,\mathrm{d}\mathrm{\Omega })`$, as described in section II-C. For the special case where the wavelet is azimuthally symmetric (i.e. invariant under azimuthal rotations), it is essentially only a function of $`\theta `$ and may be represented in terms of its Legendre expansion. In this case the harmonic representation of the wavelet coefficients is given by the product of the signal and wavelet spherical harmonic coefficients:
$$(\widehat{W_\psi ^s})_\mathrm{}m=\sqrt{\frac{4\pi }{2\mathrm{}+1}}\widehat{\psi }_\mathrm{}0^{}\widehat{s}_\mathrm{}m,$$
(17)
noting that the harmonic coefficients of an azimuthally symmetric wavelet are zero for $`m0`$. In practice, one requires that at least one of the signals, usually the wavelet, has a finite band limit so that negligible power is present in those coefficients above a certain $`\mathrm{}_{\mathrm{max}}`$. We then only need to consider $`\mathrm{}\mathrm{}_{\mathrm{max}}`$ (a detailed discussion of the determination of $`\mathrm{}_{\mathrm{max}}`$ is presented in ). Once the spherical harmonic representation of the wavelet coefficients is calculated by (17), the inverse spherical harmonic transform is applied to compute the wavelet coefficients in the Euler domain. The complexity of the fast isotropic CSWT algorithm is dominated by the spherical harmonic transforms. For a tessellation containing pixels on rings of constant latitude, a fast spherical harmonic transform may be performed (see e.g.). This reduces the complexity of the spherical harmonic transform from $`𝒪(N^4)`$ to $`𝒪(N^3)=𝒪(N_{\mathrm{pix}}^{}{}_{}{}^{3/2})`$. For certain tessellation schemes fast spherical harmonic transforms of lower complexity are also available, however these are related directly to the tessellation (e.g. ). In particular, the algorithm developed by for the ECP tessellation scales as $`𝒪(N^2\mathrm{log}N)`$.
The fast azimuthally symmetric CSWT algorithm is posed purely in harmonic space and consequently the algorithm is tessellation independent. However, we are restricted to azimuthally symmetric wavelets and lose the ability to perform the directional analysis inherent in the wavelet transform construction.
### III-E Fast directional case
We present the most general fast directional CSWT algorithm for non-azimuthally symmetric wavelets, i.e. steerable and directional wavelets, in this section. Again, the algorithm is posed purely in harmonic space and so is tessellation independent. We do, however, use the equi-angular $`_2`$ discretisation of the wavelet coefficient domain, although other discretisations may be used if FFTs are also defined on these grids. The CSWT at a particular scale (i.e. a particular $`a`$ and $`b`$) is essentially a spherical convolution, hence we apply the fast spherical convolution algorithm proposed by to evaluate the wavelet transform. The algorithm proceeds by factoring the rotation into two separate rotations, each of which involves only a constant polar rotation component. Azimuthal rotations may then be performed in harmonic space at far less computation expense than polar rotations. We subsequently rederive the fast spherical convolution algorithm developed by , as applied to our application of evaluating the CSWT. The harmonic representation of the CSWT is first presented, followed by the discretisation and fast implementation.
#### III-E1 Harmonic formulation
Substituting the spherical harmonic expansions of the wavelet and signal into the wavelet transform defined by (8), and noting the orthogonality of the spherical harmonics described by (11), yields the harmonic representation
$`W_\psi ^s(\alpha ,\beta ,\gamma )=`$ (18)
$`{\displaystyle \underset{\mathrm{}=0}{\overset{\mathrm{}_{\mathrm{max}}}{}}}{\displaystyle \underset{m=\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \underset{m^{}=\mathrm{}}{\overset{\mathrm{}}{}}}\left[D_{mm^{}}^{\mathrm{}}(\alpha ,\beta ,\gamma )\widehat{\psi }_\mathrm{}m^{}\right]^{}\widehat{s}_\mathrm{}m.`$
Again, we assume negligible power above $`\mathrm{}_{\mathrm{max}}`$ in at least one of the signals, usually the wavelet, so that the outer summation is truncated to $`\mathrm{}_{\mathrm{max}}`$. The additional summation over $`m^{}`$ and the $`D_{mm^{}}^{\mathrm{}}`$ Wigner rotation matrices that are introduced characterise the rotation of a spherical harmonic, noting that a rotated spherical harmonic may be represented by a sum of harmonics of the same $`\mathrm{}`$ :
$$\left[R(\alpha ,\beta ,\gamma )Y_\mathrm{}m\right](\omega )=\underset{m^{}=\mathrm{}}{\overset{\mathrm{}}{}}D_{m^{}m}^{\mathrm{}}(\alpha ,\beta ,\gamma )Y_\mathrm{}m^{}(\omega ).$$
The Wigner rotation matrices may be decomposed as
$$D_{mm^{}}^{\mathrm{}}(\alpha ,\beta ,\gamma )=e^{im\alpha }d_{mm^{}}^{\mathrm{}}(\beta )e^{im^{}\gamma },$$
(19)
where the real polar $`d`$-matrix is defined by
$`d_{mm^{}}^{\mathrm{}}(\beta )={\displaystyle \underset{t=\mathrm{max}(0,mm^{})}{\overset{\mathrm{min}(\mathrm{}+m,\mathrm{}m^{})}{}}}(1)^t`$
$`\times {\displaystyle \frac{\left[(\mathrm{}+m)!(\mathrm{}m)!(\mathrm{}+m^{})!(\mathrm{}m^{})!\right]^{1/2}}{(\mathrm{}+mt)!(\mathrm{}m^{}t)!(t+m^{}m)!t!}}`$
$`\times \left[\mathrm{cos}\left({\displaystyle \frac{\beta }{2}}\right)\right]^{2\mathrm{}+mm^{}2t}\left[\mathrm{sin}\left({\displaystyle \frac{\beta }{2}}\right)\right]^{m^{}m+2t},`$
and the sum over $`t`$ is defined so that the arguments of the factorials are non-negative. Recursion formulae are available to compute rapidly the Wigner rotation matrices in the basis of either complex or real spherical harmonics. We employ the recursion formulae described in in our implementation. The decomposition shown in (19) is exploited by factoring the rotation $`R(\alpha ,\beta ,\gamma )`$ into two separate rotations, both of which only contain a constant $`\pm \pi /2`$ polar rotation:
$`R(\alpha ,\beta ,\gamma )`$ $`=`$ $`R(\alpha \pi /2,\pi /2,\beta )`$ (20)
$`\times R(0,\pi /2,\gamma +\pi /2).`$
By factoring the rotation in this manner and applying the decomposition described by (19), (18) can be rewritten as
$`W_\psi ^s(\alpha ,\beta ,\gamma )=`$ (21)
$`{\displaystyle \underset{\mathrm{}=0}{\overset{\mathrm{}_{\mathrm{max}}}{}}}{\displaystyle \underset{m=\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \underset{m^{}=\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \underset{m^{\prime \prime }=\mathrm{min}(m_{\mathrm{max}},\mathrm{})}{\overset{\mathrm{min}(m_{\mathrm{max}},\mathrm{})}{}}}d_{m^{}m}^{\mathrm{}}(\pi /2)`$
$`\times d_{m^{}m^{\prime \prime }}^{\mathrm{}}(\pi /2)\widehat{\psi }_{\mathrm{}m^{\prime \prime }}^{}\widehat{s}_\mathrm{}m`$
$`\times e^{i[m(\alpha \pi /2)+m^{}\beta +m^{\prime \prime }(\gamma +\pi /2)]},`$
where the symmetry relationship $`d_{mm^{}}^{\mathrm{}}(\beta )=d_{m^{}m}^{\mathrm{}}(\beta )`$ has been applied. (A similar factoring of the rotation and harmonic space representation has been independently performed in .) Steerable wavelets may have a low effective band limit $`m_{\mathrm{max}}\mathrm{}_{\mathrm{max}}`$, in which case the the inner summation in (21) may be truncated to $`m_{\mathrm{max}}`$. For general directional wavelets this is not the case and one must use $`m_{\mathrm{max}}=\mathrm{}_{\mathrm{max}}`$.
Evaluating the harmonic formulation given by (21) directly would be no more efficient that approximating the initial integral (8) using simple quadrature. However, (21) is represented in such a way that the presence of complex exponentials may be exploited such that FFTs may be applied to evaluate rapidly the three summations simultaneously.
#### III-E2 Fast implementation
Azimuthal rotations may be applied with far less computational expense than polar rotations since they appear within complex exponentials in (21). Although the $`d`$-matrices can be evaluated reasonably quickly and reliably using recursion formulae, the basis for the fast implementation is to avoid these polar rotations as much as possible and use FFTs to evaluate rapidly all of the azimuthal rotations simultaneously. This is the motivation for factoring the rotation by (20) so that all Euler angles occur as azimuthal rotations.
The discretisation of each Euler angle may in general be arbitrary. However, to exploit standard FFT routines uniform sampling is adopted, i.e. grid $`_2`$ is used (see section III-A). As mentioned, this discretisation is inefficient since it covers the $`\mathrm{SO}(3)`$ manifold exactly twice, nevertheless it enables the use of standard FFT routines which significantly increases the speed of the algorithm. Discretising (21) in this manner and interchanging the order of summation we obtain
$`(W_\psi ^s)_{n_\alpha ,n_\beta ,n_\gamma }=`$
$`{\displaystyle \underset{m=\mathrm{}_{\mathrm{max}}}{\overset{\mathrm{}_{\mathrm{max}}}{}}}{\displaystyle \underset{m^{}=\mathrm{}_{\mathrm{max}}}{\overset{\mathrm{}_{\mathrm{max}}}{}}}{\displaystyle \underset{m^{\prime \prime }=m_{\mathrm{max}}}{\overset{m_{\mathrm{max}}}{}}}`$
$`{\displaystyle \underset{\mathrm{}=\mathrm{max}(m,m^{},m^{\prime \prime })}{\overset{\mathrm{}_{\mathrm{max}}}{}}}d_{m^{}m}^{\mathrm{}}(\pi /2)d_{m^{}m^{\prime \prime }}^{\mathrm{}}(\pi /2)\widehat{\psi }_{\mathrm{}m^{\prime \prime }}^{}\widehat{s}_\mathrm{}m`$
$`\times e^{i[m(2\pi n_\alpha /N_\alpha \pi /2)+m^{}2\pi n_\beta /N_\beta +m^{\prime \prime }(2\pi n_\gamma /N_\gamma +\pi /2)]}.`$
Shifting the indices yields
$`(W_\psi ^s)_{n_\alpha ,n_\beta ,n_\gamma }=`$ (22)
$`{\displaystyle \underset{m=0}{\overset{2\mathrm{}_{\mathrm{max}}}{}}}{\displaystyle \underset{m^{}=0}{\overset{2\mathrm{}_{\mathrm{max}}}{}}}{\displaystyle \underset{m^{\prime \prime }=0}{\overset{2m_{\mathrm{max}}}{}}}e^{i2\pi (n_\alpha m/N_\alpha +n_\beta m^{}/N_\beta +n_\gamma m^{\prime \prime }/N_\gamma )}`$
$`\times T_{m,m^{},m^{\prime \prime }},`$
where
$`T_{m,m^{},m^{\prime \prime }}=e^{i(m^{\prime \prime }m)\pi /2}`$ (23)
$`\times {\displaystyle \underset{\mathrm{}=\mathrm{max}(m,m^{},m^{\prime \prime })}{\overset{\mathrm{}_{\mathrm{max}}}{}}}d_{m^{}m}^{\mathrm{}}(\pi /2)d_{m^{}m^{\prime \prime }}^{\mathrm{}}(\pi /2)`$
$`\times \widehat{\psi }_{\mathrm{}m^{\prime \prime }}^{}\widehat{s}_\mathrm{}m`$
is extended periodically. Note that the phase shift introduced by shifting the indices of the summation in (22), shifts the $`T_{m,m^{},m^{\prime \prime }}`$ indices back. Making the associations $`N_\alpha =2\mathrm{}_{\mathrm{max}}+1`$, $`N_\beta =2\mathrm{}_{\mathrm{max}}+1`$ and $`N_\gamma =2m_{\mathrm{max}}+1`$, one notices that (22) is the unnormalised 3-dimensional inverse DFT of (23). Nyquist sampling in $`(\alpha ,\beta ,\gamma )`$ is inherently ensured by the associations made with $`\mathrm{}_{\mathrm{max}}`$ and $`m_{\mathrm{max}}`$.
The CSWT may be performed rapidly in spherical harmonic space by constructing the $`T`$-matrix of (23) from spherical harmonic coefficients and precomputed $`d`$-matrices, followed by the application of an FFT to evaluate rapidly all three Euler angles of the discretised CSWT simultaneously. The complexity of the algorithm is dominated by the computation of the $`T`$-matrix. This involves performing a 1-dimensional summation over a 3-dimensional grid, hence the algorithm is of order $`𝒪(N_\gamma N^3)`$.
Additional benefits may be realised for real signals by exploiting the resulting conjugate symmetry relationship. Memory and computational requirements may be reduced by a further factor of two for real signals by noting the conjugate symmetry relationship $`T_{m,m^{},m^{\prime \prime }}=T_{m,m^{},m^{\prime \prime }}^{}`$. In our implementation we use the complex-to-real FFT routines of the FFTW<sup>9</sup><sup>9</sup>9http://www.fftw.org/ package, which are approximately twice as fast as the equivalent complex-to-complex routines.
### III-F Comparison
We summarise the computational complexities of the various CSWT algorithms for a single pair of scales and single orientation in Table I. The complexity of each algorithm scales with the number of dilations considered and, for those algorithms that facilitate a directional analysis (i.e. all but the fast azimuthally symmetric algorithm), with the number of orientations considered. The most general fast directional algorithm provides a saving of $`𝒪(N)`$ over the direct case, where the number of pixels on the sphere and the harmonic band limit are related to $`N`$ by $`𝒪(N_{\mathrm{pix}})=𝒪(\mathrm{}_{\mathrm{max}}^{}{}_{}{}^{2})=𝒪(N^2)`$.
We implement all algorithms in Fortran 90, adopting the HEALPix tessellation of the sphere (which, incidentally, is the tessellation scheme of the WMAP CMB data ). Typical execution times for the algorithms are tabulated in Table II for a range of resolutions from 110 down to 3.4 arcminutes. The improvements provided by the fast algorithms are apparent. Indeed, it is not feasible to run the direct algorithm on data-sets with a resolution much greater than $`N_{\mathrm{pix}}5\times 10^5`$. For data-sets of practical size, such as the WMAP ($`N_{\mathrm{pix}}3\times 10^6`$) and forthcoming Planck ($`N_{\mathrm{pix}}50\times 10^6`$) CMB data, the fast implementations of the CSWT are essential.
The semi-fast algorithm is also implemented using the HEALPix tessellation. However, to perform the outer summation (integration) continual interpolation followed by down-sampling is required on each iso-latitude ring to essentially resample the data on an ECP tessellation. This increases the execution time of the implementation of the semi-fast algorithm on the HEALPix grid to an extent that the semi-fast algorithm provides little advantage over the direct algorithm. To appreciate the advantages of the semi-fast approach it must be implemented on an ECP tessellation, hence we do not tabulate the execution times for our implementation of this algorithm on the HEALPix grid as it provides an unfair comparison.
It is also important to note that although complexity scales with the number of dilations and orientations considered, execution time does not for the fast algorithms. Execution time does scale in this manner for the direct algorithm. There are a number of additional overheads associated with the fast algorithms, such as computing spherical harmonic coefficients and $`d`$-matrices. Consequently, the fast algorithms provide additional execution time savings that are not directly apparent in Table II. For example, the execution time of the fast azimuthally symmetric and directional algorithms for 10 dilations at a resolution of $`N_{\mathrm{pix}}8\times 10^5`$ ($`N_{\mathrm{side}}=256`$) are 3:08.20 and 7:06.83 (minutes:seconds) respectively, which is considerably faster than ten times the execution time of one dilation.
## IV Application
We demonstrate in this section the application of our CSWT implementation to binary Earth data. In Fig. 3 the Earth data and the corresponding spherical wavelet coefficients are shown. We use the SBW defined in section II-D to perform the analysis. This is a steerable wavelet , however our implementation is in general valid for any directional wavelet. Notice how the wavelet coefficient maps corresponding to different oriented wavelets pick out corresponding oriented structure in the data. As the dilation scale is increased, the scale of the features extracted also increases accordingly.
The ability to probe oriented structure in data defined on a sphere is of important practical use. Certain physical processes may be localised on the sphere in scale, position and orientation (e.g. signatures of cosmic strings in the CMB or correlations induced in the CMB by the nearby galaxy distribution ). Thus, analysing the statistical properties of spherical wavelet coefficients individually for a range of scales, positions and orientations may allow one to detect such effects with greater significance. Indeed, using a directional spherical wavelet analysis we have made very strong detections of non-Gaussianity in the CMB and the strongest detection made to date of the ISW effect .
## V Concluding remarks
The extension of Euclidean wavelet analysis to the sphere has been described in the framework presented by , where the stereographic projection is used to relate the spherical and Euclidean constructions. We extend the concept of the spherical dilation presented by to anisotropic dilations. Although anisotropic dilations are of practical use, the resulting basis one projects onto does not classify as a wavelet basis since perfect reconstruction is not possible.
Current and forthcoming data-sets on the sphere, of the CMB for example, are of considerable size as higher resolutions are necessary to test new physics. Consequently, we present fast algorithms to implement the CSWT as an analysis without such algorithms is not computationally feasible. A range of algorithms are described, from the direct quadrature approximation, to the semi-fast equi-angular algorithm where one rotation is performed in Fourier space, to the fast azimuthally symmetric and directional algorithms posed purely in spherical harmonic space. Posing the CSWT purely in harmonic space naturally ensures the resulting algorithms are tessellation independent. The most general fast directional algorithm provides a saving of $`𝒪(\sqrt{N_{\mathrm{pix}}})=𝒪(\mathrm{}_{\mathrm{max}})`$ over the direct implementation and may be performed down to a few arcminutes even with limited computational resources.
Data is observed on a sphere in a range of applications. In many of these cases the ability to perform a wavelet analysis would be useful. For example, spherical wavelets may be used to probe the CMB for deviations from the standard assumption of Gaussianity, or to search for compact objects embedded in the CMB, such as cosmic strings, a predicted but as yet unobserved phenomenon. The extension of wavelet analysis to the sphere enables one to probe new physics in many areas, and is facilitated on large practical data-sets by our fast directional CSWT algorithm.
## Acknowledgements
JDM is supported by a Commonwealth (Cambridge) Scholarship from the Association of Commonwealth Universities and the Cambridge Commonwealth Trust. DJM is supported by the UK Particle Physics and Astronomy Research Council (PPARC). The implementations described in this paper use the HEALPix and FFTW packages. We also acknowledge use of the YAWTb Matlab toolbox for the binary Earth data defined on the sphere. |
warning/0506/math0506308.html | ar5iv | text | # A note on: “Relaxation Oscillators with Exact Limit Cycles”.The authors are partially supported by a MCYT grant number BFM 2002-04236-C02-01. The first author is also partially supported by DURSI of Government of Catalonia “Distinció de la Generalitat de Catalunya per a la promoció de la recerca universitària”.
## 1 Introduction and statement of the main result
Our purpose in this work is to give a family of planar polynomial differential systems of the form:
$$\dot{x}=P(x,y),\dot{y}=Q(x,y),$$
(1)
for which an explicit expression of a limit cycle, that is, an isolated periodic orbit, can be given. We assume that $`P(x,y)`$ and $`Q(x,y)`$ belong to the ring of real polynomials in two variables $`[x,y]`$, and we will always assume that $`P(x,y)`$ and $`Q(x,y)`$ are coprime polynomials. We denote by $`\mathrm{d}=\mathrm{max}\{\mathrm{deg}P,\mathrm{deg}Q\}`$ and we say that $`\mathrm{d}`$ is the degree of system (1).
In the work a family of planar polynomial differential systems like (1) is studied and the existence of an explicit limit cycle is pretended to be given. The author of gives a family of systems of the form (1) with a prescribed invariant algebraic curve. This curve $`f(x,y)=0`$ has an oval surrounding the origin of coordinates. However, in there is no proof of the fact that the oval of $`f(x,y)=0`$ is an isolated periodic orbit, that is, a limit cycle. It is stated as obvious. We have been able to weaken the hypothesis appearing in , getting a bigger family of planar polynomial systems, and we have been able to show that the oval of $`f(x,y)=0`$ is a hyperbolic limit cycle.
###### Theorem 1
We consider a polynomial $`p(x)`$ such that $`p(x_e)=0`$ and $`p(x_d)=0`$ for some $`x_e<0`$ and $`x_d>0`$, $`p^{}(x_e)0`$ and $`p^{}(x_d)0`$. We assume that $`p(x)>0`$ for all $`x`$ in the interval $`(x_e,x_d)`$. We consider another polynomial $`q(x)`$ satisfying $`p(x)q(x)^21`$ for all $`x(x_e,x_d)`$ and $`q^{}(x)0`$ for all $`x(x_e,x_d)`$. Then, the algebraic curve given by $`f(x,y)=0`$ with $`f(x,y):=(yp(x)q(x))^2p(x)`$ has an oval in the band $`x_exx_d`$ which is a hyperbolic limit cycle for the following system:
$$\begin{array}{ccc}\dot{x}\hfill & =\hfill & y,\hfill \\ \dot{y}\hfill & =\hfill & \left(\frac{3}{2}q(x)p^{}(x)+p(x)q^{}(x)\right)y\frac{p^{}(x)}{2}\left(p(x)q(x)^21\right).\hfill \end{array}$$
(2)
We note that a system of the form (2) can also be viewed as an autonomous Liénard differential equation:
$$\ddot{x}+f(x)\dot{x}+g(x)=0,$$
(3)
where $`f(x)`$ and $`g(x)`$ are the polynomials given by: $`f(x):=3q(x)p^{}(x)/2p(x)q^{}(x)`$ and $`g(x):=(p^{}(x)/2)\left(p(x)q(x)^21\right)`$. Therefore, Theorem 2 may be used to construct models of Liénard differential equations exhibiting a hyperbolic limit cycle.
One of the most famous Liénard differential equation is called van der Pol equations and it appears when studying the vacuum-tube circuits. This particular equation (3) has $`f(x)=\mu (x^21)`$ and $`g(x)=1`$, with $`\mu `$, and it exhibits a unique hyperbolic limit cycle surrounding the origin. This limit cycle is shown to be non-algebraic in the work of Odani . We are not considering van der Pol’s equation since the systems described in Theorem 2 always exhibit an algebraic limit cycle. The systems given in (2) are examples of Liénard equations with an algebraic and hyperbolic limit cycle.
The same kind equations are studied in the work , but under other hypothesis for the polynomials $`p(x)`$ and $`q(x)`$, and the author of pretends to state the existence of a limit cycle. The conditions for the polynomials $`p(x)`$ and $`q(x)`$ appearing in are: $`p(x)`$ is an even polynomial, $`p(0)>0`$, there exists a value $`X>0`$ such that $`p(X)=0`$, $`p(x)<0`$ for all $`x>X`$, $`q(x)`$ is an odd polynomial and $`q(0)=0`$. All these conditions are contained in the ones that we assume in Theorem 2. However, the authors noticed that in the work the condition $`p(x)q(x)^21`$ for $`x(X,X)`$ does not appear and it is not implied by the other hypothesis. As we will see in the proof of Theorem 2, condition $`p(x)q(x)^21`$ for $`x(x_e,x_d)`$ is necessary to have a limit cycle and it cannot be avoided.
We remark that each system of the family (2) has an algebraic limit cycle, which is an oval of the real algebraic curve $`f(x,y)=0`$, and it may have other limit cycles which are not taken into consideration. These other limit cycles can be contained in $`f(x,y)=0`$ or not. If they are contained in an invariant curve, we can treat them with the same methods described in this note. For instance, it can be shown that the system (2) with $`p(x)=(1x^2)(4x^2)(9x^2)`$ and $`q(x)=x/100`$ has $`3`$ hyperbolic limit cycles all of them contained in the corresponding invariant algebraic curve $`f(x,y)=0`$.
In order to prove Theorem 2, we need some results relating ovals of curves of the form $`f(x,y)=0`$ with the fact that the oval of such a curve is a limit cycle. These preliminary results are stated in Section 2. Once we have stated these previous results, we prove Theorem 2 in Section 3.
## 2 Preliminary results
We are considering limit cycles which are contained in a real curve $`f(x,y)=0`$, which does not need to be algebraic. This fact leads us to the definition of invariant of a system (1).
###### Definition 2
Let us consider an open set $`𝒰^2`$ and a $`𝒞^1(𝒰)`$ real function denoted by $`f(x,y):𝒰^2`$. We say that $`f(x,y)`$ is an invariant for a system (1) if
$$P(x,y)\frac{f}{x}(x,y)+Q(x,y)\frac{f}{y}(x,y)=k(x,y)f(x,y),$$
(4)
with $`k(x,y)`$ a polynomial of degree lower or equal than $`\mathrm{d}1`$, where d is the degree of the system. This polynomial $`k(x,y)`$ is called the cofactor of $`f(x,y)`$.
In case that $`f(x,y)`$ is a polynomial we say that $`f(x,y)=0`$ is an invariant algebraic curve for system (1). We notice that if $`f(x,y)`$ is an invariant of system (1) and $`f(x,y)=0`$ defines a curve in the real plane, then the function $`P(x,y)(f/x)+Q(x,y)(f/y)`$ equals zero on the points such that $`f(x,y)=0`$. This fact implies that the real curve $`f(x,y)=0`$ is formed by orbits of system (1). In particular if $`f(x,y)=0`$ contains an oval without any singular point of system (1), this oval is a periodic orbit of system (1).
As well as invariant curves, the other objects taken into consideration in this paper are limit cycles. A limit cycle of system (1) is an isolated periodic orbit. Let $`\gamma `$ be a limit cycle for system (1). We say that $`\gamma `$ is stable if there exists a neighborhood such that all the orbits starting in it have $`\gamma `$ as $`\omega `$–limit set. We say that $`\gamma `$ is unstable if there is a neighborhood such that all the orbits starting in it have $`\gamma `$ as $`\alpha `$-limit set. There might be limit cycles which are neither stable nor unstable. These limit cycles have a neighborhood such that in the interior of the limit cycle all the orbits have $`\gamma `$ as $`\omega `$-limit set and in the exterior of $`\gamma `$ all the orbits have $`\gamma `$ as $`\alpha `$-limit set. Or the other way round: the orbits of the interior have $`\gamma `$ as $`\alpha `$-limit set and the orbits in the exterior have $`\gamma `$ as $`\omega `$-limit set. In this case, we say that $`\gamma `$ is semi-stable. Any limit cycle $`\gamma `$ of a system (1) is either stable, unstable or semi-stable as it is stated in .
A classical known result, given in the book of Perko , let us distinguish the hyperbolicity of a limit cycle. If we consider $`\gamma (t)`$ a periodic orbit of system (1) of period $`T`$, we may compute the finite value given by the following integral $`_0^T\mathrm{div}(\gamma (t))𝑑t`$, where $`\mathrm{div}(x,y)=(P/x)+(Q/y)`$ is called the divergence of system (1). It can be shown that if $`_0^T\mathrm{div}(\gamma (t))𝑑t<0`$, then $`\gamma `$ is a stable limit cycle, if $`_0^T\mathrm{div}(\gamma (t))𝑑t>0`$, then $`\gamma `$ is a unstable limit cycle and if $`_0^T\mathrm{div}(\gamma (t))𝑑t=0`$, then $`\gamma `$ may be a stable, unstable or semi-stable limit cycle or it may belong to a continuous band of cycles. When the quantity $`_0^T\mathrm{div}(\gamma (t))𝑑t`$ is different from zero, we say that the limit cycle $`\gamma `$ is hyperbolic. We notice that if $`_0^T\mathrm{div}(\gamma (t))𝑑t0`$, then the periodic orbit $`\gamma `$ is a limit cycle (either stable or unstable). We are going to use this property to ensure that a periodic orbit is a limit cycle, that is, that it does not belong to a continuous band of cycles.
We relate limit cycles with invariants in the following way. We assume that we have a periodic orbit $`\gamma `$ of system (1) which is given in an implicit way, that is, there exists an invariant curve $`f(x,y)=0`$ such that $`\gamma \{(x,y)|f(x,y)=0\}`$. In order to have a smooth curve $`f(x,y)=0`$ defining the periodic orbit, we will assume that $`f(p)0`$ for any $`p\gamma `$, that is, the gradient vector of $`f(x,y)`$ is different from zero in all the points of $`\gamma `$. Then we have the following result stated and proved in .
###### Theorem 3
Let us consider a system (1) and $`\gamma (t)`$ a periodic orbit of period $`T>0`$. Assume that $`f:𝒰^2`$ is an invariant curve with $`\gamma \{(x,y)|f(x,y)=0\}`$ and let $`k(x,y)`$ be the cofactor of $`f(x,y)`$ as given in (4). We assume that $`f(p)0`$ for any $`p\gamma `$. Then,
$$_0^Tk(\gamma (t))𝑑t=_0^T\mathrm{div}(\gamma (t))𝑑t.$$
(5)
Hence, we have an alternative way to compute the value $`_0^T\mathrm{div}(\gamma (t))𝑑t`$. In the family of planar polynomial differential systems which we are considering, that is, the one described in Theorem 2, we will not be able to directly compute the value $`_0^T\mathrm{div}(\gamma (t))𝑑t`$. This is due to the fact that we are not considering a fixed system with a concrete periodic orbit, but a family of systems each one with a different periodic orbit and, thus, the expression of the integrand is too general to be manipulated. By Theorem 5, we can also compute the value $`_0^Tk(\gamma (t))𝑑t`$ but this integral is as much difficult as the previous one. That’s why we are going to use the fact that for any $`w`$:
$$_0^T\mathrm{div}(\gamma (t))𝑑t=_0^T\mathrm{div}(\gamma (t))𝑑t+w\left(_0^T\mathrm{div}(\gamma (t))𝑑t_0^Tk(\gamma (t))𝑑t\right).$$
The integrand in the right hand side of this equality will be chosen strictly positive or negative in all the interval of integration for a suitable value of $`w`$. Therefore, the value of the integral will be different from zero. Using these steps, we will be able to prove that the oval of the invariant curve described in Theorem 2 is a limit cycle of the corresponding system.
## 3 Proof of Theorem 2
In order to prove this theorem, we first show that $`f(x,y)=0`$, where $`f(x,y):=(yp(x)q(x))^2p(x)`$, is an invariant algebraic curve of system (2). Straightforward computations show that:
$$\begin{array}{c}y\left(\frac{f}{x}\right)+\left[\left(\frac{3}{2}q(x)p^{}(x)+p(x)q^{}(x)\right)y\frac{p^{}(x)}{2}\left(p(x)q(x)^21\right)\right]\left(\frac{f}{y}\right)=\hfill \\ =q(x)p^{}(x)f(x,y),\hfill \end{array}$$
and, thus, we have that $`f(x,y)=0`$ is an invariant algebraic curve for system (2) with cofactor $`k(x,y):=q(x)p^{}(x)`$.
Since $`p(x_e)=p(x_d)=0`$ for the values $`x_e<0`$ and $`x_d>0`$ and $`p(x)>0`$ for $`x(x_e,x_d)`$, we deduce that $`f(x,y)=0`$ has an oval in the band $`x_exx_d`$ surrounding the origin of coordinates, which can be parameterized in two parts by:
$$x(\tau )=\tau ,y_\pm (\tau )=p(\tau )q(\tau )\pm \sqrt{p(\tau )},$$
(6)
with $`\tau (x_e,x_d)`$. We are going to prove that this oval does not contain any singular point of system (2), and then, we will have that it defines a periodic orbit of the system. The singular points of system (2) have coordinates of the form $`(a,0)`$ where the value $`a`$ is a root of the polynomial $`p^{}(x)(p(x)q(x)^21)`$. We have that $`f(a,0)=p(a)(p(a)q(a)^21)`$ and $`(p(a)q(a)^21)`$ is different from zero in all the closed interval $`a[x_e,x_d]`$ by the hypothesis that $`p(x)q(x)^21`$ for all $`x(x_e,x_d)`$ and $`p(x_e)=p(x_d)=0`$.
Here we notice that the assumption $`p(x)q(x)^21`$ for $`x(x_e,x_d)`$ is necessary for the oval of $`f(x,y)=0`$ to be a limit cycle. In , this assumption is not given. We notice that an oval of an invariant algebraic curve of a system may contain singular points of the system, and in such a case, it is not even a periodic orbit.
Since $`x_e`$ and $`x_d`$ are simple zeroes of $`p(x)`$ and $`p(x)>0`$ for all $`x`$ in the interval $`(x_e,x_d)`$ we have that there is no singular point of system (2) on the oval given by $`f(x,y)=0`$ and parameterized by (6). From this fact and that $`f(x,y)=0`$ is an invariant algebraic curve of the system, we deduce that this oval is a periodic orbit of system (1). We denote this periodic orbit by $`\gamma `$ for the rest of the proof. We note that we do not know the parameterization of $`\gamma `$ as explicit solution of system (2), that is, we do not know the periodic function $`\gamma (t):=(\gamma _1(t),\gamma _2(t))`$ such that $`d\gamma _1(t)/dt=\gamma _2(t)`$ and
$`{\displaystyle \frac{d\gamma _2(t)}{dt}}`$ $`=`$ $`\left({\displaystyle \frac{3}{2}}q(\gamma _1(t))p^{}(\gamma _1(t))+p(\gamma _1(t))q^{}(\gamma _1(t))\right)\gamma _2(t)`$
$`{\displaystyle \frac{p^{}(\gamma _1(t))}{2}}\left(p(\gamma _1(t))q(\gamma _1(t))^21\right),`$
for all $`t`$. We do neither know its period $`T>0`$ but we have been able to show its existence by using the invariant algebraic curve $`f(x,y)=0`$ and its properties in relation with system (2).
Finally, we need to prove that the periodic orbit $`\gamma `$ is a hyperbolic limit cycle. To do so, we are going to show that the value of the integral $`_0^T\mathrm{div}(\gamma (t))𝑑t`$ is different from zero. Since we do not know $`\gamma (t)`$ nor the period $`T`$, we use the parameterization of the oval $`\gamma `$ given in (6). In order to get the correct sign of the integral $`_0^T\mathrm{div}(\gamma (t))𝑑t`$, we need to know the sense of the flow over $`\gamma `$. We take the point of coordinates $`(x_d,0)`$, which belongs to $`\gamma `$, and we have that the vector field defined by system (2) on that point is $`(0,p^{}(x_d)/2)`$ because $`p(x_d)=0`$. Since $`p(x_d)=0`$, $`p^{}(x_d)0`$ and $`p(x)>0`$ in the interval $`(x_e,x_d)`$, we deduce that $`p^{}(x_d)<0`$. Hence, the sense of the flow over $`\gamma `$ is clockwise. We can write the following equality, using the parameterization (6):
$`{\displaystyle _0^T}\mathrm{div}(\gamma (t))𝑑t`$ $`=`$ $`{\displaystyle _{x_e}^{x_d}}{\displaystyle \frac{\mathrm{div}(x(\tau ),y_+(\tau ))}{y_+(\tau )}}𝑑\tau +{\displaystyle _{x_d}^{x_e}}{\displaystyle \frac{\mathrm{div}(x(\tau ),y_{}(\tau ))}{y_{}(\tau )}}𝑑\tau `$
$`=`$ $`{\displaystyle _{x_e}^{x_d}}\left({\displaystyle \frac{\mathrm{div}(x(\tau ),y_+(\tau ))}{y_+(\tau )}}{\displaystyle \frac{\mathrm{div}(x(\tau ),y_{}(\tau ))}{y_{}(\tau )}}\right)𝑑\tau .`$
We note that the divergence of system (2) is $`\mathrm{div}(x,y)=3q(x)p^{}(x)/2+p(x)q^{}(x)`$, and substituting this expression in the former equality, we get:
$$_0^T\mathrm{div}(\gamma (t))𝑑t=_{x_e}^{x_d}\frac{3q(\tau )p^{}(\tau )2p(\tau )q^{}(\tau )}{(p(\tau )q(\tau )^21)\sqrt{p(\tau )}}𝑑\tau .$$
This integral is well defined because we are assuming that $`(p(\tau )q(\tau )^21)p(\tau )`$ is different from zero for $`\tau (x_e,x_d)`$. However, we are not able to distinguish if its value is positive, negative or zero. By using the same reasonings, we can write the following equality:
$$_0^Tk(\gamma (t))𝑑t=_{x_e}^{x_d}\frac{2q(\tau )p^{}(\tau )}{(p(\tau )q(\tau )^21)\sqrt{p(\tau )}}𝑑\tau .$$
Using Theorem 5, we have that, for any value of $`w`$:
$`{\displaystyle _0^T}\mathrm{div}(\gamma (t))𝑑t`$ $`=`$ $`{\displaystyle _0^T}\mathrm{div}(\gamma (t))𝑑t+w\left({\displaystyle _0^T}\mathrm{div}(\gamma (t))𝑑t{\displaystyle _0^T}k(\gamma (t))𝑑t\right)`$
$`=`$ $`{\displaystyle _{x_e}^{x_d}}{\displaystyle \frac{(w+3)q(\tau )p^{}(\tau )2(1+w)p(\tau )q^{}(\tau )}{(p(\tau )q(\tau )^21)\sqrt{p(\tau )}}}𝑑\tau .`$
Taking $`w=3`$, we get:
$$_0^T\mathrm{div}(\gamma (t))𝑑t=_{x_e}^{x_d}\frac{4p(\tau )q^{}(\tau )}{(p(\tau )q(\tau )^21)\sqrt{p(\tau )}}𝑑\tau $$
(7)
The hypothesis of Theorem 2 on the sign of the polynomials $`p(x)`$ and $`q(x)`$ in the interval $`x(x_e,x_d)`$ are $`p(x)>0`$, $`p(x)q(x)^21`$ and $`q^{}(x)0`$. Therefore the integrand of the right hand side of (7) is strictly positive or negative in all the interval $`\tau (x_e,x_d)`$. We deduce that the value of the integral cannot be zero and, hence, the periodic orbit $`\gamma `$ is a hyperbolic limit cycle as we wanted to show.
We would also like to characterize if this limit cycle is stable or unstable, so we are going to study the sign of the integrand in the right hand side of (7). Since $`p(x)q(x)^21`$ for $`x(x_e,x_d)`$ and $`p(x_d)q(x_d)^21=1`$ (because $`p(x_d)=0`$), we deduce that $`p(x)q(x)^21<0`$ for $`x(x_e,x_d)`$. Therefore, we have that the integrand in the right hand side of (7) is strictly positive if $`q^{}(\tau )<0`$ for all $`\tau (x_e,x_d)`$ and strictly negative if $`q^{}(\tau )>0`$ for all $`\tau (x_e,x_d)`$. We can state that the hyperbolic limit cycle $`\gamma `$ is stable if $`q^{}(0)>0`$ and unstable if $`q^{}(0)<0`$.
We also note that in the work , the expression of an example of a more general limit cycle for a family of planar polynomial differential systems is pretended to be given. In fact, we are going to show that in the case that the oval of this example is a limit cycle, we are in the same family of systems as written in (2), that is, the one described in Theorem 2.
In the following planar polynomial differential system is given as an example of a more general family of systems with an explicit limit cycle.
$$\begin{array}{ccc}\dot{x}\hfill & =\hfill & y,\hfill \\ \dot{y}\hfill & =\hfill & \{p^{}(x)[(m+r)h(x)p(x)^{r1}+(m+1)q(x)]+h^{}(x)p(x)^r\hfill \\ & & +p(x)q^{}(x)\}ymp(x)p^{}(x)([h(x)p(x)^{r1}+q(x)]^2p(x)^{2m2}),\hfill \end{array}$$
(8)
where $`p(x)`$ is an even polynomial, $`q(x)`$ and $`h(x)`$ are odd polynomials, $`m=n+1/2`$, $`n`$ is an integer number with $`n0`$ and $`r`$ is an integer number with $`r2`$. Moreover, it is assumed that $`p(0)>0`$ and there exists a value $`X>0`$ such that $`p(X)=0`$, $`p(x)<0`$ for all $`x>X`$ and $`q(0)=0`$.
Some straightforward computations show that system (8) exhibits the invariant algebraic curve $`f(x,y)=0`$ with $`f(x,y):=(yh(x)p(x)^rq(x)p(x))^2p(x)^{2n+1}`$ and with cofactor $`k(x,y):=(2n+1)p^{}(x)\left[h(x)p(x)^{r1}+q(x)\right]`$. We have that $`f(x,y)=0`$ has an oval in the band $`XxX`$ which can be parameterized by:
$$x(\tau )=\tau ,y(\tau )=h(\tau )p(\tau )^r+p(\tau )q(\tau )\pm p(\tau )^n\sqrt{p(\tau )}.$$
In order to show that this oval is a periodic orbit of system (8), we only need to show that it does not contain any singular point of the system. The singular points of system (8) in the band $`|x|X`$ are of the form $`(a,0)`$ where $`a`$ is a root of the polynomial $`p(x)p^{}(x)\left(\left[h(x)p(x)^{r1}+q(x)\right]^2p(x)^{2n1}\right)`$, because $`m=n+1/2`$. We notice that, unless $`n=0`$, the points with coordinates $`(X,0)`$ and $`(X,0)`$ are singular points of the system which are contained in the oval of $`f(x,y)=0`$. Therefore, if $`n>0`$, we have that the oval of $`f(x,y)=0`$ cannot be a limit cycle. If $`n=0`$, we can consider the polynomial $`\stackrel{~}{q}(x):=q(x)+h(x)p(x)^{r1}`$ and we have that system (8) coincides with system (2) with polynomials $`p(x)`$ and $`\stackrel{~}{q}(x)`$. Therefore, this is not an example of a more general limit cycle. |
warning/0506/math0506216.html | ar5iv | text | # Minimal Volume Entropy on Graphs
## 1.Introduction
Let $`(X,g)`$ be a compact connected Riemannian manifold of nonpositive curvature. It was shown by A. Manning \[Man\] that the topological entropy $`h_{top}(g)`$ of the geodesic flow is equal to the volume entropy $`h_{vol}(g)`$ of the manifold
$$h_{vol}(g)=\underset{r\mathrm{}}{lim}\frac{1}{r}\mathrm{log}(vol(B(x,r))),$$
where $`B(x,r)`$ is the ball of radius $`r`$ centered at some point $`x`$ in the universal cover $`\stackrel{~}{X}`$ of $`X`$. G. Besson, G. Courtois and S. Gallot \[BCG\] proved that if $`X`$ has dimension at least three and carries a rank one locally symmetric metric $`g_0`$, then for every Riemannian metric $`g`$ such that $`vol(X,g)=vol(X,g_0)`$, the inequality
$$h_{vol}(g)h_{vol}(g_0)$$
holds, with equality if and only if $`g`$ is isometric to $`g_0`$. This solved a conjecture mainly due to M. Gromov \[Gro\], which had been proved earlier by A. Katok \[Kat\] for metrics in the conformal class of the hyperbolic metric on a compact orientable surface.
In this paper, we are interested in the analogous problem for finite graphs, endowed with metrics obtained by varying the length $`\mathrm{}(e)`$ of the edges $`eEX`$ of a graph $`X`$. Regular and biregular trees, as rank one buildings (see \[BT\]), are non-archimedian analogs of rank one symmetric spaces, and they carry many lattices (see \[BL\] for instance). As in the case of Riemannian manifolds, it is well known that the volume entropy of a finite metric graph is equal to the topological entropy of the geodesic flow on its universal cover (see \[Gui\]) and also equal to the critical exponent of its fundamental group acting on its universal covering tree (\[Bou\]).
We show that for any graph, among all (normalized) metrics on it, there exists a unique metric minimizing the volume entropy. We further give explicit formulas for the minimal volume entropy and the metric realizing it.
Theorem. Let $`X`$ be a finite connected graph such that the valency at each vertex $`x`$, which we denote by $`k_x+1`$, is at least $`3`$. Then there is a unique normalized length distance minimizing the volume entropy $`h_{\mathrm{vol}}(d)`$. The minimal volume entropy is
$$h_{\mathrm{min}}=\frac{1}{2}\underset{xVX}{}(k_x+1)\mathrm{log}k_x,$$
and the entropy minimizing length distance $`d=d_{\mathrm{}}`$ is given by
$$eEX,\mathrm{}(e)=\frac{\mathrm{log}(k_{i(e)}k_{t(e)})}{\underset{xVX}{}(k_x+1)\mathrm{log}k_x}.$$
In the special case of a regular graph, which is the analog of a Riemannian manifold carrying a locally symmetric metric, the volume entropy is minimized by the metric for which all the edges have the same length. This special case was independently shown by I. Kapovich and T. Nagnibeda (\[KN\]).
The above theorem has an analog for graphs of groups (see \[Ser\]) as well, as in Proposition 9. Finally, using Proposition 9, we show that for a $`n`$-sheeted covering graph of groups (see for instance \[Bas\]) $`\varphi :(Y,H_{})(X,G_{})`$, there holds, for the proper definition of the volume of a metric graph of groups (see section 4),
$$h_{vol}(Y,H_{},d)vol(Y,H_{},d)nh_{vol}(X,G_{},d_0)vol(X,G_{},d_0),$$
and that the equality holds if and only if the length distance $`d`$ on $`(Y,H_{})`$ is an entropy-minimizing length distance among the length distances of the same volume, and the map $`\varphi `$ is a metric covering from $`(Y,H_{},d)`$ to $`(X,G_{},\lambda d_0),`$ for some $`\lambda >0`$. This can be considered as an analog of the main theorem in \[BCG\] for graphs.
Acknowledgement. We are grateful to F. Paulin for suggesting this problem as well as for helpful discussions and encouragement. We would also like to thank H. Furstenberg and G. Margulis for insightful discussions and H. Furstenberg for explaining us the relevant work \[FK\]. We thank I. Kapovich for pointing out the preprint of I. Rivin (\[Riv\]).
## 2. Volume entropy and path growth
Let us consider a nonempty connected unoriented finite graph $`X`$ without any terminal vertex. We will denote the set of vertices by $`VX`$ and the set of oriented edges of $`X`$ by $`EX`$. We denote again by $`X`$ the geometric realization of $`X`$. For every edge $`e`$, let us denote by $`i(e)`$ and $`t(e)`$ the initial and the terminal vertex of $`e`$, respectively. We define a length distance $`d`$ on $`X`$ by assigning a positive real number $`\mathrm{}(e)=\mathrm{}(\overline{e})`$ for each unoriented edge $`\{e,\overline{e}\}`$ of $`X`$, and by letting $`d=d_{\mathrm{}}:X\times X[0,\mathrm{}[`$ be the maximal distance which makes each half-edge of an edge $`e`$ containing a vertex, isometric to $`[0,\frac{\mathrm{}(e)}{2}]`$. For a length distance $`d_{\mathrm{}}`$, let $`l_{\mathrm{max}}=\underset{eEX}{\mathrm{max}}\mathrm{}(e)`$ and $`l_{\mathrm{min}}=\underset{eEX}{\mathrm{min}}\mathrm{}(e)`$. Define the volume of $`X`$ by
$$vol(X,d)=\frac{1}{2}\underset{eEX}{}\mathrm{}(e),$$
i.e., the sum of lengths of the unoriented edges. We denote by $`\mathrm{\Delta }(X)`$ the set of all length distances $`d=d_{\mathrm{}}`$ on $`X`$ normalized so that $`vol(X,d)=1.`$
For a fixed length distance $`d`$, let us consider a universal covering tree $`\stackrel{~}{X}X`$ equipped with the lifted distance $`\stackrel{~}{d}`$ of $`d`$. For any connected subset $`S`$ of $`\stackrel{~}{X}`$, let us denote by $`\mathrm{}(S)`$ the sum of the lengths of (the maximal pieces of) the edges in $`S`$. We define the volume entropy $`h_{vol}(d)=h_{\mathrm{vol}}(X,d)`$ as
$$h_{vol}(d)=\underset{r\mathrm{}}{lim\; sup}\frac{1}{r}\mathrm{log}\mathrm{}(B(x_0,r)),$$
where $`B(x_0,r)=B_d(x_0,r)`$ is the ball of radius $`r`$ with center a fixed vertex $`x_0`$ in $`(\stackrel{~}{X},\stackrel{~}{d})`$. The entropy $`h_{vol}(d)`$ does not depend on the base point $`x_0`$, and we may sum either on the oriented or on the non-oriented edges. Note also the homogeneity property
$$h_{vol}(d_\alpha \mathrm{})=\frac{1}{\alpha }h_{vol}(d_{\mathrm{}}),()$$
for every $`\alpha >0`$. Remark that $`h_{\mathrm{vol}}(X,d)\mathrm{vol}(X,d)`$ is invariant under dilations, therefore to minimize the entropy with constant volume, it suffices to consider the length metrics of volume $`1`$.
If $`\pi _1X`$ is not cyclic, or equivalently if $`X`$ has no terminal vertices and is not reduced to one cycle, then $`h_{\mathrm{vol}}=h_{\mathrm{vol}}(d)`$ is strictly positive, which we will assume from now on (see for instance \[Bou\]). It was shown by Roblin (\[Robl\]) that the upper limit above is in fact a limit. This implies that as $`r\mathrm{}`$,
$$\mathrm{}(B(x_0,r))=e^{h_{\mathrm{vol}(d)}r+o(r)}.$$
By a metric path of length $`r`$ in $`X`$, we mean the image of a local isometry $`f:[0,r]X`$. Note that the endpoint of a metric path is not necessarily a vertex. By a combinatorial $`n`$-path of length $`r`$ in $`X`$, we mean a path $`\underset{¯}{p}=e_1e_2\mathrm{}e_n`$ of consecutive edges in $`X`$ without backtracking such that $`_{j=1}^{n1}\mathrm{}(e_j)<r_{j=1}^n\mathrm{}(e_j)`$. A combinatorial path is a combinatorial $`n`$-path for some $`n`$.
###### Lemma 1.
Let $`N_r(x_0)`$ be the cardinality of the set of combinatorial paths of length $`r`$ in $`\stackrel{~}{X}`$ starting at $`x_0V\stackrel{~}{X}`$. Then the number $`N_r(x_0)`$ satisfies
$$\underset{r\mathrm{}}{lim\; sup}\frac{\mathrm{log}N_r(x_0)}{r}=\underset{r\mathrm{}}{lim}\frac{\mathrm{log}N_r(x_0)}{r}=h_{vol}.$$
Proof. It follows directly from $`\mathrm{}(B(x_0,r))=e^{(h_{\mathrm{vol}}+o(1))r}`$ that for any $`l>0`$,
$$\underset{r\mathrm{}}{lim\; sup}\frac{\mathrm{log}\mathrm{}(B(x_0,r)B(x_0,rl))}{r}=\underset{r\mathrm{}}{lim}\frac{\mathrm{log}\mathrm{}(B(x_0,r)B(x_0,rl))}{r}=h_{vol}.$$
Now let $`N_r^{}(x_0)`$ be the cardinality of the set of metric paths of length $`r`$ starting at $`x_0`$. As $`\stackrel{~}{X}`$ has no terminal vertices, for any $`l>0`$,
$$lN_{rl}^{}(x_0)\mathrm{}(B(x_0,r)B(x_0,rl))lN_r^{}(x_0).$$
Therefore
$$\underset{r\mathrm{}}{lim\; sup}\frac{\mathrm{log}N_r^{}(x_0)}{r}=\underset{r\mathrm{}}{lim}\frac{\mathrm{log}N_r^{}(x_0)}{r}=h_{vol}.$$
It is clear that we get a combinatorial path of length $`r`$ by continuing a metric path of length $`r`$ until it meets a vertex. Also, two distinct combinatorial paths of length $`r`$ cannot be extensions of one metric path of length $`r`$ by the strict inequality in the definition of a combinatorial path. It follows that $`N_r(x_0)=N_r^{}(x_0),`$ thus $`N_r(x_0)`$ has the same exponential growth rate as $`N_r^{}(x_0)`$, which is $`h_{vol}`$. ∎
Let $`A=A(X)`$ be the edge adjacency matrix of $`X`$, i.e. a $`|EX|\times |EX|`$ matrix defined by $`A_{ef}=\rho _{ef}`$, where $`\rho _{ef}`$ has value $`1`$ if $`ef`$ is a combinatorial path, i.e. if $`t(e)=i(f)`$ and $`\overline{e}f`$, and value 0 otherwise. It is easy to see that the entry $`A_{ef}^n`$ is nonzero if and only if there is a combinatorial $`n`$-path starting with $`e`$ and ending with $`f`$. (Note that the definition of $`\rho _{ef}`$ implies that such a path does not have backtracking.)
Let us show that for any connected graph without any terminal vertex, which is not a cycle, the matrix $`A`$ is irreducible. Recall that a matrix $`M`$ is reducible if there exists a permutation matrix $`P`$ such that $`PMP^1`$ is a block diagonal matrix of at least two nontrivial blocks. A matrix is irreducible if it is not reducible. It is clear that a nonnegative matrix $`M`$ is reducible if and only if $`M+M^t`$ is reducible. The matrix $`A+A^t`$ has entries $`b_{ef}`$ with $`b_{ef}`$ nonzero if either $`ef`$ or $`fe`$ is a combinatorial path. Let $``$ be the equivalence relation on $`EX`$ generated by the relation $`ef`$ if $`\rho _{ef}=1`$ (i.e if $`ef`$ is a path without backtracking).
###### Lemma 2.
If $`EX/`$ consists of only one element, then $`A`$ is irreducible.
Proof. By contradiction. Assume that $`A`$ is reducible, and hence $`A+A^t`$ is also reducible. This means that it is possible to find two complementary subsets $`U`$ and $`V`$ in $`EX`$ such that $`b_{uv}=0`$ for any $`uU`$ and $`vV`$. But then clearly no element $`u`$ of $`U`$ can be equivalent to an element $`v`$ of $`V`$, and in particular there are at least two equivalence classes for the equivalence relation $``$. ∎
###### Proposition 3.
Let $`X`$ be a connected graph without any terminal vertex. Then the matrix $`A`$ is irreducible if and only if $`X`$ has a vertex of valency at least three.
Proof. We first claim that for any $`e,fEX`$, either $`ef`$ or $`e\overline{f}`$. Let $`𝐞`$ and $`𝐟`$ be the unoriented edges underlying $`e`$ and $`f`$. There is a shortest (thus without backtracking) unoriented path in $`X`$ linking $`𝐞`$ and $`𝐟`$. Thus, there is an combinatorial (oriented) path in $`X`$ linking either $`e`$ or $`\overline{e}`$ to either $`f`$ or $`\overline{f}`$. If $`\underset{¯}{p}=e_1\mathrm{}e_s`$ is a combinatorial path, then so is $`\overline{\underset{¯}{p}}=\overline{e_s}\mathrm{}\overline{e_1}`$. We deduce that, for any two edges $`e_1`$ and $`e_2`$, $`e_1e_2\overline{e_2}\overline{e_1}`$. Applying this to $`e,\overline{e}`$ and $`f,\overline{f}`$ yields the claim.
Now if $`X`$ has only bivalent vertices (i.e. if $`X`$ is a cycle, as $`X`$ has no terminal vertices) then it is easy to see that $`A`$ is not irreducible. Conversely, let $`X`$ be connected and let $`x`$ be a vertex of valence at least three, with outgoing edges $`e_1,e_2`$ and $`e_3`$ (these may be loops). Then $`\overline{e_1}e_3`$ and $`\overline{e_2}e_3`$ hence $`\overline{e}_1\overline{e_2}`$. But $`\overline{e_1}e_2`$, and therefore $`e_2\overline{e_2}`$. Now for any edge $`e`$ in $`EX`$, either $`ee_2`$ or $`e\overline{e_2}`$ by the claim above. In both cases, we have $`ee_2`$. This implies that $`EX/`$ has only one element and $`A`$ is irreducible. ∎
Now consider the matrix $`A^{}=A^{}(d,h)`$ defined by $`A_{ef}^{}=\rho _{ef}e^{h\mathrm{}(f)}`$, depending on $`h`$ and the length distance $`d_{\mathrm{}}`$ on $`X`$. The matrix $`A^{}`$ is clearly irreducible since $`A`$ is irreducible.
###### Theorem 4.
Let $`X`$ be a connected finite graph without any terminal vertex, which is not a cycle, endowed with a length distance $`d=d_{\mathrm{}}`$. The volume entropy $`h_{\mathrm{vol}}`$ is the only positive constant $`h`$ such that the following system of linear equations with unknowns $`(x_e)_{eEX}`$ has a solution with $`x_e>0`$ for every $`eEX`$.
$$x_e=\underset{fEX}{}\rho _{ef}e^{h\mathrm{}(f)}x_f,\mathrm{𝑓𝑜𝑟}\mathrm{𝑎𝑙𝑙}eEX.()$$
Proof. By the assumption on the graph, for every $`h>0`$, we can apply Perron-Frobenius theorem (see \[Gan\] for example) to the irreducible nonnegative matrix $`A^{}=(\rho _{ef}e^{h\mathrm{}(f)})`$, which says that the spectral radius of the matrix $`A^{}(h)`$ is a positive eigenvalue $`\lambda (h)`$, which is simple, with an eigenvector $`(x_e=x_e(h))`$ whose entries are all positive. The function $`\lambda :_0_0`$ is clearly a continuous function of $`h`$ since the characteristic function of the matrix $`A^{}`$ is a polynomial of $`\{e^{h\mathrm{}(e)}:eEX\}`$, and $`\lambda (0)1`$ since $`\lambda (0)`$ is the spectral radius of an irreducible nonzero matrix $`A^{}(0)`$ of nonnegative integer coefficients. Also, $`\lambda (h)0`$ as $`h\mathrm{}`$, since the coefficients of $`A^{}(h)`$ tends to $`0`$ as $`h\mathrm{}`$. By the mean value theorem, there exists an $`h`$ satisfying $`\lambda (h)=1`$.
Now assume that $`h>0`$ satisfies $`()`$ for some positive $`x_e`$’s. Fix an arbitrary edge $`eEX`$ (any edge $`e`$ satisfies $`x_e>0`$ by Perron-Frobenius theorem), and choose a vertex $`x_0`$ in $`\stackrel{~}{X}`$ which is an initial vertex of a fixed lift $`\stackrel{~}{e}`$ of $`e`$ in $`\stackrel{~}{X}`$. Let us fix a positive constant $`r`$.
Let $`P_r(e)`$ be the set of combinatorial paths of length $`r`$ in $`X`$ starting with $`e`$. We will denote a combinatorial path in $`X`$ by $`\underset{¯}{p}=e_1e_2\mathrm{}e_n`$, its terminal edge by $`t(\underset{¯}{p})=e_n`$ and its metric length by $`\mathrm{}(\underset{¯}{p})=_{i=1}^n\mathrm{}(e_i)`$. Denote by $`𝒫_n(e)`$ (resp. $`𝒫_n^{}(e)`$) the set of combinatorial $`k`$-paths of length $`r`$ with $`kn`$ (resp. combinatorial $`n`$-paths of length strictly less than $`r`$) in $`X`$ starting with $`e`$. Remark that $`𝒫_n(e)𝒫_n^{}(e)=\mathrm{}`$ and if $`n`$ is large enough, $`𝒫_n(e)=P_r(e)`$ and $`𝒫_n^{}(e)=\mathrm{}`$.
Let us rewrite the equation (\**) as
$$e^{h\mathrm{}(e)}x_e=\underset{\underset{¯}{p}𝒫_2(e)𝒫_2^{}(e)}{}e^{h\mathrm{}(\underset{¯}{p})}x_{t(\underset{¯}{p})}.$$
Let us replace each $`x_{t(\underset{¯}{p})}`$ in the above equation by $`\underset{fEX}{}\rho _{t(\underset{¯}{p})f}e^{h\mathrm{}(f)}x_f`$ whenever $`\mathrm{}(\underset{¯}{p})<r`$, i.e. when $`\underset{¯}{p}𝒫_2^{}(e)`$. The resulting equation is
$$e^{h\mathrm{}(e)}x_e=\underset{\underset{¯}{p}𝒫_3(e)𝒫_3^{}(e)}{}e^{h\mathrm{}(\underset{¯}{p})}x_{t(\underset{¯}{p})}.$$
Repeat this process: at each step, for each $`\underset{¯}{p}𝒫_n^{}(e)`$, replace $`x_{t(\underset{¯}{p})}`$ on the right hand side of the previous equation by $`\underset{fEX}{}\rho _{t(\underset{¯}{p})f}e^{h\mathrm{}(f)}x_f`$, to get
$$e^{h\mathrm{}(e)}x_e=\underset{\underset{¯}{p}𝒫_{n+1}(e)𝒫_{n+1}^{}(e)}{}e^{h\mathrm{}(\underset{¯}{p})x_{t(\underset{¯}{p})}}.$$
For $`n`$ large enough, the resulting equation is
$$e^{h\mathrm{}(e)}x_e=\underset{\underset{¯}{p}P_r(e)}{}e^{h\mathrm{}(\underset{¯}{p})}x_{t(\underset{¯}{p})}.$$
(In the case when the lengths of the edges are all equal to $`1`$ and $`r`$ is a positive integer, we continue until we get the equation $`\underset{¯}{x}=A^{r1}\underset{¯}{x}`$.)
Then in the resulting equation, the number of times each $`x_f`$ appears on the right hand side is exactly the number $`N_r(e,f)`$ of combinatorial paths of length $`r`$ in $`\stackrel{~}{X}`$ with initial edge $`\stackrel{~}{e}`$ and terminal edge some lift of $`f`$ in $`\stackrel{~}{X}`$. Note also that the metric length of such a path is at least $`r`$ and less than $`r+l_{\mathrm{max}}`$. Thus
$$\underset{fEX}{}N_r(e,f)e^{h(r+l_{\mathrm{max}})}x_fe^{h\mathrm{}(e)}x_e\underset{fEX}{}N_r(e,f)e^{hr}x_f.$$
Now if $`h`$ is strictly greater than the volume entropy $`h_{\mathrm{vol}}`$, then
$`0<e^{h\mathrm{}(e)}x_e`$ $`{\displaystyle N_r(e,f)e^{hr}x_f}N_r(x_0)e^{hr}{\displaystyle x_f}`$
$`e^{(h_{\mathrm{vol}}+0(1))r}e^{hr}{\displaystyle x_f}=e^{r(h_{\mathrm{vol}}h+o(1))}{\displaystyle x_f}0`$
as $`r`$ goes to infinity, which is a contradiction. On the other hand, suppose that $`h`$ is strictly smaller than the volume entropy $`h_{\mathrm{vol}}`$. As $`N_r(x_0)=\underset{e,fEX,i(\stackrel{~}{e})=x_0}{}N_r(e,f)`$, there exist some $`e`$ and $`f`$, depending on $`r`$, such that $`N_r(e,f)e^{hr}\frac{1}{|EX|^2}N_r(x_0)e^{hr}`$. Since
$`e^{h\mathrm{}(e)}x_e`$ $`N_r(e,f)e^{h(r+l_{\mathrm{max}})}x_f{\displaystyle \frac{1}{|EX|^2}}e^{(h_{\mathrm{vol}}+o(1))r}e^{h(r+l_{\mathrm{max}})}x_f`$
$`{\displaystyle \frac{1}{|EX|^2}}e^{(h_{\mathrm{vol}}h+o(1))r}e^{hl_{\mathrm{max}}}\underset{fEX}{\mathrm{min}}\{x_f\},`$
and the latter goes to infinity as $`r`$ goes to infinity, it follows that $`x_e=\mathrm{}`$, which is again a contradiction. We conclude that $`h`$ is equal to the volume entropy $`h_{\mathrm{vol}}`$. ∎
Remark. Hersonsky and Hubbard showed in \[HH\] that the Hausdorff dimension of the limit set of a Schottky subgroup of the automorphism group of a simplicial tree satisfies similar systems of equations.
## 3. Minimal volume entropy
In this section, we prove the main theorem announced in the introduction, using Theorem 4.
###### Theorem 5.
Let $`X`$ be a finite connected graph such that the valency at each vertex $`x`$, which we denote by $`k_x+1`$, is at least $`3`$. Then there is a unique $`d`$ in $`\mathrm{\Delta }(X)`$ minimizing the volume entropy $`h_{\mathrm{vol}}(d)`$. The minimal volume entropy is
$$h_{\mathrm{min}}(X)=\frac{1}{2}\underset{xVX}{}(k_x+1)\mathrm{log}k_x,$$
and the entropy minimizing length distance $`d=d_{\mathrm{}}`$ is characterized by
$$\mathrm{}(e)=\frac{\mathrm{log}k_{i(e)}k_{t(e)}}{\underset{xVX}{}(k_x+1)\mathrm{log}k_x},eEX.$$
Remark. Since we can eliminate all the vertices of valency two without changing the entropy, the existence of $`d`$ in $`\mathrm{\Delta }(X)`$ minimizing the volume entropy, with minimal value given by the same formula, holds for any graph who does not have a terminal vertex and is not isometric to a circle. What is uniquely defined at such a minimum is the length of each connected component of $`X`$ where the vertices of valency at least three are removed.
Proof. By assumption, $`k_x2`$ for every $`xVX`$. By Theorem 5, the volume entropy $`h=h_{\mathrm{vol}}`$ satisfies
$$x_e=\underset{fEX}{}\rho _{ef}e^{h\mathrm{}(f)}x_f,$$
(1)
for each edge $`eEX`$ for some positive $`x_e`$’s. Set $`y_e=e^{h\mathrm{}(e)}x_e>0`$ for each edge $`e`$. Then the above equations implies
$$e^{h\mathrm{}(e)}y_e=\underset{fEX}{}\rho _{ef}y_fk_{t(e)}\underset{fEX,\rho _{ef}=1}{}y_f^{1/k_{i(f)}}.$$
The last inequality is simply the inequality between the arithmetic mean and the geometric mean of $`y_f`$’s, since there are exactly $`k_{t(e)}=k_{i(f)}`$ edges $`f`$ such that $`\rho _{ef}=1`$. Multiplying over all the edges, we get
$$\underset{eEX}{}e^{h\mathrm{}_e}y_e\underset{eEX}{}(k_{t(e)}\underset{fEX,\rho _{ef}=1}{}y_f^{1/k_{i(f)}}).$$
On the right hand side of the equation, each term $`y_f^{1/k_{i(f)}}`$ appears exactly $`k_{i(f)}`$ times, since each edge $`f`$ follows exactly $`k_{i(f)}`$ edges with terminal vertex $`i(f)`$. Canceling $`\underset{eEX}{}y_e>0`$ from each side, we get
$$e^{2h}\underset{eEX}{}k_{t(e)}=\underset{xVX}{}k_x^{(k_x+1)},$$
(2)
since $`\underset{eEX}{}\mathrm{}(e)=2`$. The equality holds if and only if equality in the inequality (1) holds for each $`eEX`$, i.e. the $`y_f`$’s, for $`fEX`$ following $`e`$, are all equal.
Suppose that the equality in the inequality (2) holds. In particular,
$$h=\frac{1}{2}\underset{xVX}{}(k_x+1)\mathrm{log}k_x.$$
Since the valency at each vertex is at least $`3`$, we can choose another edge $`gf`$ followed by $`e`$ and conclude that $`y_f`$ depends only on the initial vertex $`i(f)`$ of $`f`$. Let $`z_{i(f)}=y_f>0`$. Then the equation $`()`$ in Theorem 4 amounts to
$$e^{h\mathrm{}(e)}z_{i(e)}=\underset{fEX}{}\rho _{ef}z_{i(f)}=k_{t(e)}z_{t(e)}.$$
Since $`\mathrm{}(e)=\mathrm{}(\overline{e})`$, we also have $`e^{h\mathrm{}(e)}z_{t(e)}=k_{i(e)}z_{i(e)}`$. Thus $`z_{i(e)}/z_{t(e)}=k_{t(e)}/e^{h\mathrm{}(e)}=e^{h\mathrm{}(e)}/k_{i(e)}`$ and
$$e^{h\mathrm{}(e)}=\sqrt{k_{i(e)}k_{t(e)}},$$
so that
$$\mathrm{}(e)=\frac{\mathrm{log}k_{i(e)}k_{t(e)}}{\underset{xVX}{}(k_x+1)\mathrm{log}k_x}.$$
(3)
In particular, $`\mathrm{}`$ is uniquely defined by this formula. The length distance defined by the formula (3) clearly satisfies the equations $`()`$, with
$$h=\frac{1}{2}\underset{xVX}{}(k_x+1)\mathrm{log}k_x,$$
and $`x_e`$’s defined, uniquely up to constant, by setting
$$\frac{e^{h\mathrm{}(e)}x_e}{e^{h\mathrm{}(f)}x_f}=\sqrt{\frac{k_{t(e)}}{k_{i(e)}}},$$
for every $`f`$ such that $`i(f)=t(e)`$. It is clearly well-defined, since if there is a cycle consisting of consecutive edges $`(e_1,e_2,\mathrm{},e_n,e_{n+1}=e_1)`$, then
$`y_{e_n}=y_{e_{n1}}\sqrt{{\displaystyle \frac{k_{i(e_{n1})}}{k_{i(e_n)}}}}=\mathrm{}=y_{e_1}{\displaystyle \underset{j=2}{\overset{j=n}{}}}\sqrt{{\displaystyle \frac{k_{i(e_{j1})}}{k_{i(e_j)}}}}=y_{e_1}\sqrt{{\displaystyle \frac{k_{i(e_1)}}{k_{i(e_n)}}}}.`$
By uniqueness in Theorem 4, the positive number $`h`$ given above is the volume entropy of the given length distance, and it is the minimal entropy of the graph.∎
###### Corollary 6.
If $`X`$ is a $`(k_1+1,k_2+1)`$-biregular graph, with $`k_1>1,k_2>1`$, then the volume entropy of the normalized length distances on $`X`$ is minimized exactly when the lengths of the edges are all equal, and the minimal volume entropy is $`\frac{|EX|}{4}\mathrm{log}(k_1k_2)`$.
Proof. Suppose that $`X`$ a $`(k_1+1,k_2+1)`$-biregular graph, i.e. $`k_{i(e)}k_{t(e)}=k_1k_2`$ for any edge $`e`$. Let $`d=d_{\mathrm{}}\mathrm{\Delta }(X)`$ be the entropy-minimizing length distance. Then $`\mathrm{}(e)=\frac{1}{2h}\mathrm{log}(k_1k_2)`$ does not depend on $`e`$, thus $`\mathrm{}(e)=\frac{2}{|EX|}`$. From $`e^{h\mathrm{}(e)}=\sqrt{k_{i(e)}k_{t(e)}},`$ the volume entropy of this length distance is $`h=\frac{|EX|}{4}\mathrm{log}(k_1k_2)`$. ∎
###### Corollary 7.
If $`X`$ is a $`(k+1)`$-regular graph, with $`k>1`$, then the volume entropy of the normalized length distances on $`X`$ is minimized exactly when the lengths of the edges are all equal, and the minimal volume entropy is $`\frac{|EX|}{2}\mathrm{log}k`$.
Proof. This is a special case of the above corollary with $`k_1=k_2=k`$.∎
Remark. The last corollary appears implicitly in a preprint of I. Rivin (\[Riv\]). There he considers graphs with weights given on the vertices rather than the edges. The directed line graph $`L(X)`$ of a graph $`X`$ is an oriented graph defined so that $`VL(X)=EX`$ and $`EL(X)=\{(a,b)EX^2:t(a)=i(b),a\overline{b}\}`$. To a given set of weights on the edges $`\{\mathrm{}(e)\}_{EX}`$, is associated a set of weights $`\{\mathrm{}^{}(x)\}_{VL(X)}`$ on the vertices of $`L(X)`$. One can sees that paths on $`X`$ without backtracking correspond to paths with backtracking on $`L(X)`$, see \[Riv\] page 14. The minimum of volume entropy of the graph $`L(X)`$ with vertex weights $`h((\mathrm{}^{}(x)))_{VL(X)}`$ (computed by I. Rivin) lies in the image of the map $`(\mathrm{}(e))(\mathrm{}^{}(x))`$ only when the graph is regular. It seems that for general graphs, one result cannot be deduced from the other.
Remark. Corollary 7 was also shown independently by I. Kapovich and T. Nagnibeda \[KN\] by a different method (using random walks). Note that one of their main results, on the minimal entropy among all graphs having a fixed fundamental group, can be deduced from Theorem 5 as in the following corollary. A special case when the graph has a highly transitive automorphism group had been shown earlier by G. Robert (\[Rob\]).
###### Corollary 8.
(\[KN\] Theorem B) Consider the set of all finite metric graphs without a vertex of valency one or two, whose fundamental group is a free group of given rank $`r2`$. Then among volume 1 length metrics, the volume entropy is minimized by any (regular) trivalent graph in this set, with the metric assigning the same length for every edge.
Proof. Let $`(X,d)`$ be such a graph. Suppose that there is a vertex $`x`$ of valency $`k_x+1`$ strictly greater than three, with outgoing edges $`e_1,\mathrm{},e_{k_x+1}`$. Let us introduce a new vertex $`y`$ and a new edge $`f`$, and replace $`x`$ and its outgoing edges $`e_1,\mathrm{},e_{k_x+1}`$, by two vertices $`x`$ and $`y`$, with outgoing edges $`f,e_3,\mathrm{},e_{k_x+1}`$ and $`e_1,e_2,\overline{f}`$, respectively. Repeat the operation on $`x`$, until the valency of $`x`$ reduces to three, to get a new graph $`X^{}`$. The graph $`X^{}`$ has $`k_x2`$ more vertices than $`X`$, all with valency three.
Let $`d_0`$ and $`d_0^{}`$ be the unique normalized entropy-minimizing length distances on $`X`$ and $`X^{}`$, respectively. By the formula in Theorem 5, since for $`t3`$, $`(t+1)\mathrm{log}t>(t1)3\mathrm{log}2`$, it follows that
$`h_{\mathrm{vol}}(X,d)`$ $`h_{\mathrm{vol}}(X,d_0)={\displaystyle \frac{1}{2}}\underset{zVX\{x\}}{{\displaystyle }}(k_z+1)\mathrm{log}k_z+(k_x+1)\mathrm{log}k_x`$
$`>{\displaystyle \frac{1}{2}}\underset{zVX\{x\}}{{\displaystyle }}(k_z+1)\mathrm{log}k_z+(k_x1)3\mathrm{log}2=h_{\mathrm{vol}}(X^{},d_0^{}).`$
Repeat the operation until we get a regular trivalent graph. Now by Corollary 7, the volume entropy is minimized when all the edges have the same length. ∎
## 4. Entropy for graphs of groups
As another corollary of Theorem 5, let us show the analogous result of Theorem 5 for graphs of groups. Let $`(X,G_{})`$ be any finite connected graph of finite groups. (Basic references for graphs of groups are \[Ser\] and \[Bas\].) Let $`T`$ be a (Bass-Serre) universal covering tree of $`(X,G_{})`$ and let $`p:TX`$ be the canonical projection. The degree of a vertex $`x`$ of $`(X,G_{})`$ is defined by
$$\underset{eEX,i(e)=x}{}\frac{|G_x|}{|G_e|}.$$
It is easy to see that it is equal to the valency of any lift of $`x`$ in $`VT`$, and we will denote it again by $`k_x+1`$. We define a length distance $`d_{\mathrm{}}`$ on $`(X,G_{})`$ as a length distance $`d_{\mathrm{}}`$ on the underlying graph $`X`$. The volume of $`(X,G_{},d_{\mathrm{}})`$ for a given length distance $`d_{\mathrm{}}`$ on $`(X,G_{})`$, is defined by
$$vol_{\mathrm{}}(X,G_{})=\frac{1}{2}\underset{eEX}{}\frac{\mathrm{}(e)}{|G_e|}.$$
Note that in the case where $`\mathrm{}(e)`$ is equal to $`1`$ for every edge $`e`$ and $`T`$ is $`k`$-regular, the volume $`vol_{\mathrm{}}(X,G_{})`$ is $`k/2`$ times the usual definition of the volume $`_{xVX}1/|G_x|`$ of a graph of groups since $`k=_{eEX,i(e)=x}|G_x|/|G_e|`$. The volume entropy $`h_{\mathrm{vol}}(X,G_{},d_{\mathrm{}})`$ of $`(X,G_{},d_{\mathrm{}})`$ is defined to be the exponential growth of the balls in $`T`$ for the lifted metric as in the case of graphs.
###### Proposition 9.
Let $`(X,G_{})`$ be a finite graph of finite groups such that the degree at each vertex $`x`$ of $`(X,G_{})`$ is at least three. Among the normalized (i.e. volume one) length distances on $`(X,G_{})`$, there exists a unique normalized length distance minimizing the volume entropy. At this minimum, the length of each edge is proportional to $`\mathrm{log}(k_{i(e)}k_{t(e)})`$ and the minimal volume entropy is
$$h_{\mathrm{min}}(X,G_{})=\frac{1}{2}\underset{xVX}{}\frac{(k_x+1)\mathrm{log}k_x}{|G_x|}.$$
Proof. Let $`\mathrm{\Gamma }`$ be a fundamental group of the graph of groups $`(X,G_{})`$. There exists a free normal subgroup $`\mathrm{\Gamma }^{}`$ of $`\mathrm{\Gamma }`$ of finite index (see \[Ser\]), say $`m`$. The group $`\mathrm{\Gamma }^{}`$ acts freely on $`T`$, hence the quotient graph $`X^{}=\mathrm{\Gamma }^{}\backslash T`$ is a finite connected graph. It is easy to see that each $`x`$ in $`VX`$ (resp. $`e`$ in $`EX`$) has $`\frac{m}{|G_x|}`$(resp. $`\frac{m}{|G_e|}`$) lifts in $`VX^{}`$ (resp. $`EX^{}`$) by the canonical map $`\pi :X^{}X`$, since
$$m=[\mathrm{\Gamma }:\mathrm{\Gamma }^{}]=\frac{_{x^{}VX^{}}1}{_{xVX}1/|G_x|}$$
(see \[Bas\] for example). It is clear that $`\mathrm{}^{}(e)=\mathrm{}(\pi (e))`$ and the valency $`k_y+1`$ is equal to the degree $`k_{\pi (y)}+1`$. Any length distance $`d_{\mathrm{}}`$ of volume one on $`(X,G_{})`$ can be lifted to $`X^{}`$ to define a length distance $`d_{\mathrm{}^{}}`$ normalized so that
$$vol_{\mathrm{}^{}}(X^{})=\frac{1}{2}\underset{eEX^{}}{}\mathrm{}^{}(e)=\frac{1}{2}\underset{eEX}{}\frac{m}{|G_e|}\mathrm{}(e)=m.$$
The volume entropy of $`(X^{},d_{\mathrm{}}^{})`$ is equal to the volume entropy of $`(X,G_{},d_{\mathrm{}})`$ as they have the same universal covering metric tree. By the homogeneity property $`()`$, we can apply Theorem 5 to conclude that among the length distances of volume $`m`$ on $`X^{}`$, there exists a unique entropy-minimizing length distance $`d_0^{}=d_{\mathrm{}^{}}`$ on $`X^{}`$. By uniqueness in Theorem 5, the length distance $`d_0^{}`$ is invariant under the group $`\mathrm{\Gamma }/\mathrm{\Gamma }^{}`$. In particular, there is a normalized length distance $`d_0=d_{\mathrm{}}`$ on $`(X,G_{})`$ whose lift to $`X^{}`$ defines $`d_0^{}`$. The minimal volume entropy of $`(X,G_{})`$ is clearly the volume entropy of $`(X^{},d_0^{})`$ since for any length distance $`d`$ on $`(X,G_{})`$,
$$h_{\mathrm{vol}}(X,G_{},d)=h(X^{},d^{})h(X^{},d_0^{})=h_{\mathrm{vol}}(X,G_{},d_0),$$
where $`d^{}`$ is defined by the lift of $`d`$ on $`X^{}`$. Since the length $`\mathrm{}^{}(e)`$ of an edge $`e`$ is proportional to $`\mathrm{log}(k_{i(e)}k_{t(e)})=\mathrm{log}(k_{\pi (i(e))}k_{\pi (t(e))})`$ for every edge $`e`$ in $`EX^{}`$, so is true for every edge $`e`$ in $`EX`$. Since each vertex $`x`$ in $`VX`$ appears $`\frac{m}{|G_x|}`$ times in $`X^{}`$ and the degree $`k_x+1`$ is equal to the valency $`k_x^{}+1`$ of any lift $`x^{}\pi ^1(x)`$ of $`x`$ in $`X^{}`$, the minimal volume entropy of $`(X,G_{})`$ is
$`h_{d_0}(X,G_{})`$ $`=h(X^{},d_0^{})={\displaystyle \frac{1}{m}}h(X^{},{\displaystyle \frac{1}{m}}d_0^{})={\displaystyle \frac{1}{2m}}\underset{x^{}VX^{}}{{\displaystyle }}(k_x^{}+1)\mathrm{log}k_x^{}`$
$`={\displaystyle \frac{1}{2m}}\underset{xVX}{{\displaystyle }}{\displaystyle \frac{m}{|G_x|}}(k_x+1)\mathrm{log}k_x={\displaystyle \frac{1}{2}}\underset{xVX}{{\displaystyle }}{\displaystyle \frac{(k_x+1)\mathrm{log}k_x}{|G_x|}}.`$
Now we want to consider a more general situation than in Proposition 9. The main theorem in \[BCG\] says that if $`f:(Y,g)(X,g_0)`$ is a continuous map of non-zero degree between compact connected $`n`$-dimensional Riemannian manifolds and $`g_0`$ is a locally symmetric metric with negative curvature, then
$$h^n(Y,g)\mathrm{vol}(Y,g)|\mathrm{deg}f|h^n(X,g_0)\mathrm{vol}(X,g_0),$$
and the equality holds if and only if $`f`$ is homotopic to a Riemannian covering.
Let $`(X,G_{},d_0=d_{\mathrm{}})`$ be a finite (connected) graph of finite groups endowed with the normalized length distance minimizing the volume entropy. Let $`(Y,H_{},d)`$ be another finite graph of finite groups with a length distance. Let $`\varphi =(\varphi ,\varphi _{},\gamma _{}):(Y,H_{})(X,G_{})`$ be a (Bass-Serre) covering of graphs of groups (see \[Bas\]). The value
$$n:=\underset{y\varphi ^1(x)}{}\frac{|G_x|}{|H_y|}=\underset{f\varphi ^1(e)}{}\frac{|G_e|}{|H_f|}$$
does not depend on the vertex $`x`$ nor on the edge $`e`$ of $`X`$ since the graph $`X`$ is connected, and it is an integer. A covering graph of groups with the above $`n`$ is said to be $`n`$-sheeted (see \[Lim\]).
When $`(Y,H_{})`$ and $`(X,G_{})`$ are graphs (of trivial groups), the next corollary can be considered as an analog of the main theorem in \[BCG\].
###### Corollary 10.
Let $`\varphi :(Y,H_{})(X,G_{})`$ be a $`n`$-sheeted covering of graphs of groups and let $`d_0`$ be the entropy-minimizing length distance on $`(X,G_{})`$ of volume one. Suppose that the degree at each vertex of $`(X,G_{})`$ and $`(Y,H_{})`$ is at least three. Then there holds
$$h_{vol}(Y,H_{},d)vol(Y,H_{},d)nh_{vol}(X,G_{},d_0)vol(X,G_{},d_0).$$
The equality holds if and only if the length distance $`d`$ on $`(Y,H_{})`$ is a length distance minimizing entropy among the length distances of the same volume, and in that case the map $`\varphi `$ is a metric covering from $`(Y,H_{},d)`$ to $`(X,G_{},\lambda d_0),`$ for some $`\lambda >0`$.
Proof. By the homogeneity property $`()`$, we may assume that $`\mathrm{vol}(Y,H_{},d)=1`$. Applying Proposition 9 to $`(Y,H_{})`$ and $`(X,G_{})`$, it follows that there exists a unique length distance $`d_0^{}=d_{\mathrm{}^{}}`$ on $`Y`$ minimizing the volume entropy and that
$`h_{\mathrm{vol}}(Y,H_{},d)`$ $`h_{\mathrm{min}}(Y,H_{})={\displaystyle \frac{1}{2}}\underset{yVY}{{\displaystyle }}{\displaystyle \frac{(k_y+1)\mathrm{log}k_y}{|H_y|}}={\displaystyle \frac{1}{2}}\underset{xVX}{{\displaystyle }}\underset{y\varphi ^1(x)}{{\displaystyle }}{\displaystyle \frac{(k_x+1)\mathrm{log}k_x}{|H_y|}}`$
$`={\displaystyle \frac{1}{2}}n\underset{xVX}{{\displaystyle }}{\displaystyle \frac{(k_x+1)\mathrm{log}k_x}{|G_x|}}=nh_{\mathrm{min}}(X,G_{})=nh_{\mathrm{vol}}(X,G_{},d_0).`$
By Proposition 9, the equality holds if and only if $`d=d_0^{}`$. In that case, the length of each edge $`e`$ in $`EY`$ is proportional to $`\mathrm{log}(k_{i(e)}k_{t(e)})=\mathrm{log}(k_{i(\varphi (e))}k_{t(\varphi (e))})`$, thus proportional to the length of the edge $`\varphi (e)`$. More precisely, let $`\mathrm{}^{}(e)=c^{}\mathrm{log}(k_{i(e)}k_{t(e)})`$ for every $`eEY`$ and let $`\mathrm{}(e)=c\mathrm{log}(k_{i(e)}k_{t(e)})`$ for every $`eEX`$. From the assumption $`\mathrm{vol}_{\mathrm{}}(X,G_{})=\mathrm{vol}_{\mathrm{}^{}}(Y,H_{})=1`$, it follows that
$`1`$ $`={\displaystyle \frac{1}{2}}\underset{gEY}{{\displaystyle }}{\displaystyle \frac{c^{}\mathrm{log}(k_{i(g)}k_{t(g)})}{|H_g|}}={\displaystyle \frac{1}{2}}\underset{eEX}{{\displaystyle }}\underset{g\varphi ^1(e)}{{\displaystyle }}{\displaystyle \frac{c^{}\mathrm{log}(k_{i(g)}k_{t(g)})}{|H_g|}}={\displaystyle \frac{1}{2}}\underset{eEX}{{\displaystyle }}{\displaystyle \frac{nc^{}\mathrm{log}(k_{i(e)}k_{t(e)})}{|G_e|}},`$
and therefore
$$c^{}=\frac{1}{\frac{n}{2}\underset{eEX}{}\frac{\mathrm{log}(k_{i(e)}k_{t(e)})}{|G_e|}}=\frac{c}{n},$$
in other words, $`\mathrm{}^{}(e)=\mathrm{}(e)/n.`$
We conclude that for any length distance $`d`$ on $`(Y,H_{})`$, there holds
$$h_{vol}(Y,H_{},d)vol(Y,H_{},d)nh_{vol}(X,G_{})vol(X,G_{},d_0).$$
By Proposition 9 the equality holds if and only if $`d`$ is proportional to $`d_0^{}`$, say $`d=\lambda nd_0^{}`$ for some $`\lambda >0`$. Then the length of each edge $`e`$ in $`(Y,H_{},d)`$ is $`\lambda \mathrm{}(\varphi (e))`$, and the map $`\varphi `$ is a metric covering from $`(Y,d_0^{})`$ to $`(X,\lambda d_0)`$. ∎
##
Yale University, New Haven, CT 06520-8283, USA and
ENS-Paris, UMR 8553 CNRS, 45 rue d’Ulm, 75230 Paris Cedex 05, France
seonhee.lim@yale.edu, Seonhee.Lim@ens.fr |
warning/0506/cond-mat0506374.html | ar5iv | text | # Phase diagram of the excitonic insulator
## Abstract
Motivated by recent experiments, which give strong evidence for an excitonic insulating phase in $`\mathrm{TmSe}_{0.45}\mathrm{Te}_{0.55}`$, we developed a scheme to quantitatively construct, for generic two-band models, the phase diagram of an excitonic insulator. As a first application of our approach, we calculated the phase diagram for an effective mass two-band model with long-range Coulomb interaction. The shielded potential approximation is used to derive a generalized gap equation controlling for positive (negative) energy gaps the transition from a semi-conducting (semi-metallic) phase to an insulating phase. Numerical results, obtained within the quasi-static approximation, show a steeple-like phase diagram in contrast to long-standing expectations.
, , ,
The possibility of an excitonic insulator (EI) phase, separating, below a critical temperature, a semiconducting from a semi-metallic phase, has been predicted by theorists more than three decades ago . However, experimental efforts to establish this phase in actual materials largely failed. It is only until recently, that detailed experimental investigations of $`\mathrm{TmSe}_{0.45}\mathrm{Te}_{0.55}`$ suggested the existence of an EI phase in this compound . The pressure dependence of the electrical resistivity below $`270\mathrm{K}`$, for instance, strongly points towards an emerging EI phase . Further evidence for collective behavior which may have its origin in an EI phase comes from the linear increase of the thermal conductance and diffusivity at very low temperatures .
Under the assumption that the external pressure controls the energy gap $`\mathrm{E}_\mathrm{g}`$, the resistivity data have been used to construct a phase diagram for $`\mathrm{TmSe}_{0.45}\mathrm{Te}_{0.55}`$ in the $`\mathrm{E}_\mathrm{g}`$-T plane . Although experimental data strongly suggest that this phase diagram is the phase diagram of an EI, to unambiguously decide if this interpretation is correct requires further theoretical examination, taking the relevant parts of the electronic structure of the material into account. However, even for the simplest two-band models, a quantitative phase diagram for an EI has never been calculated. As a first step towards a theoretical scrutiny of the phases of the $`\mathrm{Tm}[\mathrm{Se},\mathrm{Te}]`$ system it is therefore appropriate to present here such a calculation.
In close analogy to the strong-coupling theory of superconductivity, we employed a matrix propagator formalism. Within a two-band model, the anomalous or off-diagonal (in the band indices $`i=1,2`$) self-energy $`\mathrm{\Sigma }_{12}(𝐤,i\omega _n)`$ describing the pairing between conduction and valence band electrons serves as an order parameter: $`\mathrm{\Sigma }_{12}(𝐤,i\omega _n)0`$ signals the existence of the EI phase. Our selfconsistent approximation enables us to take a variety of physical processes into account and results in a nonlinear functional equation for $`\mathrm{\Sigma }_{12}(𝐤,i\omega _n)`$. Linearizing this equation in the vicinity of the phase boundary, where $`\mathrm{\Sigma }_{12}(𝐤,i\omega _n)`$ is small, yields a generalized “gap equation”. The phase boundary $`\mathrm{T}_\mathrm{c}(\mathrm{E}_g)`$ can then be found by mapping out the T-$`\mathrm{E}_g`$ range for which the “gap equation” has nontrivial solutions.
We applied this scheme to an isotropic, effective mass two-band model for valence and conduction band electrons interacting via the long-range Coulomb potential $`V_0(𝐪)=4\pi e^2/\epsilon _0q^2`$. The energy gap $`\mathrm{E}_\mathrm{g}`$ is indirect ($`\mathrm{\Gamma }`$-X) and can be positive or negative. Within the shielded potential approximation the generalized gap equation for the real part of the interband self-energy reads
$`\mathrm{\Delta }(𝐤,\stackrel{~}{e}_i)`$ $`=`$ $`{\displaystyle }{\displaystyle \frac{d𝐤^{}}{(2\pi )^3}}[V(𝐤𝐤^{},\stackrel{~}{e}_i,\stackrel{~}{e}_1^{})B_2(\stackrel{~}{e}_2^{},\stackrel{~}{e}_1^{})\mathrm{\Delta }(𝐤^{},\stackrel{~}{e}_1^{})`$ (1)
$`+`$ $`V(𝐤𝐤^{},\stackrel{~}{e}_i,\stackrel{~}{e}_2^{})B_1(\stackrel{~}{e}_1^{},\stackrel{~}{e}_2^{})\mathrm{\Delta }(𝐤^{},\stackrel{~}{e}_2^{})],`$
where $`\stackrel{~}{e}_i(𝐤)=e_i(𝐤)\mu _i`$, $`B_i(x,y)=(xy)/[(xy)^2+\gamma _i(𝐤)^\mathrm{𝟐}]`$, and $`V(𝐤𝐤^{},\stackrel{~}{e}_i,\stackrel{~}{e}_j^{})=\mathrm{Re}V_s^r(𝐤𝐤^{},\stackrel{~}{e}_i\stackrel{~}{e}_j^{})n_F(\stackrel{~}{e}_j^{})\mathrm{P}(d\omega /\pi )1/(\stackrel{~}{e}_i\stackrel{~}{e}_j^{}\omega )\mathrm{Im}V_s^r(𝐤𝐤^{},\stackrel{~}{e}_i\stackrel{~}{e}_j^{})n_B(\omega )`$ with $`V_s^r(𝐪)=V_0(𝐪)/\epsilon ^r(𝐪,\omega )`$ the dynamically screened Coulomb potential. To derive Eq. (1) we employed a quasi-particle approximation for the intraband propagators with renormalized band dispersions and lifetimes given by $`e_i(𝐤)=ϵ_i(𝐤)+\mathrm{Re}\mathrm{\Sigma }_{ii}^r(𝐤,e_i(𝐤)\mu _i)`$ and $`\gamma _i(𝐤)=\mathrm{Im}\mathrm{\Sigma }_{ii}^r(𝐤,e_i(𝐤)\mu _i)`$, respectively. The chemical potentials $`\mu _i`$ are measured from the respective band extrema.
The full analysis of Eq. (1) is the subject of a forthcoming publication . Here, we focus on the quasi-static approximation, $`V_s^r(𝐤𝐤^{},\stackrel{~}{e}_i\stackrel{~}{e}_j^{})V_s^r(𝐤𝐤^{},0)`$, which simplifies the gap equation enormously. Except for very small band overlaps (very small Fermi surfaces), we expect this approximation to work reasonably well, as it does for intraband self-energies. For equal band masses and temperature independent screening, the quasi-static approximation reduces Eq. (1) to
$`\mathrm{\Delta }(u)`$ $`=`$ $`{\displaystyle _{u_1}^{\mathrm{}}}𝑑u^{}V(u,u^{}){\displaystyle \frac{\mathrm{tanh}u^{}}{2u^{}}}\mathrm{\Delta }(u^{})`$ (2)
$`V(u,u^{})`$ $`=`$ $`\sqrt{{\displaystyle \frac{1}{4\pi ^2\mathrm{k}_\mathrm{B}\mathrm{T}}}}{\displaystyle \frac{1}{\mathrm{k}}}\mathrm{log}\left[{\displaystyle \frac{(\mathrm{k}+\mathrm{k}^{})^2+\kappa ^2}{(\mathrm{k}\mathrm{k}^{})^2+\kappa ^2}}\right]`$ (3)
with $`\mathrm{k}=\sqrt{uu_1}`$, $`\mathrm{k}^{}=\sqrt{u^{}u_1}`$, $`u_1=\mathrm{E}_\mathrm{g}/4\mathrm{k}_\mathrm{B}\mathrm{T}`$, and $`\kappa ^2=(2\sqrt{|\mathrm{E}_\mathrm{g}|}/\pi \mathrm{k}_\mathrm{B}\mathrm{T})\theta (\mathrm{E}_\mathrm{g})`$. To construct the phase boundary $`\mathrm{T}_\mathrm{c}(\mathrm{E}_g`$), we discretize Eq. (2) and determine, for fixed $`\mathrm{E}_\mathrm{g}`$, the temperature $`\mathrm{T}=\mathrm{T}_\mathrm{c}`$ for which the determinant of the coefficient matrix of the resulting system of linear equations vanishes. For $`\mathrm{E}_\mathrm{g}<0`$ this approach can be directly applied, whereas for $`\mathrm{E}_\mathrm{g}>0`$, the logarithmic singularity of the kernel has to be removed first .
The phase boundary $`\mathrm{T}_\mathrm{c}(\mathrm{E}_g`$) is presented in Fig. 1, measuring energy and temperatures in units of the exciton Rydberg $`R_0`$. Above $`T_10.45`$, the EI phase is unstable. Below $`T_1`$, we find a steeple-like phase boundary which strongly discriminates between $`\mathrm{E}_\mathrm{g}>0`$ and $`\mathrm{E}_\mathrm{g}<0`$. For $`\mathrm{E}_\mathrm{g}>0`$, $`T_c(\mathrm{E}_g`$) smoothly decreases to zero at $`\mathrm{E}_\mathrm{g}=1`$, the critical band gap, above which the EI phase cannot exist. For $`\mathrm{E}_\mathrm{g}<0`$, in contrast, $`T_c(\mathrm{E}_g)`$ initially drops extremely fast, within a few percent of $`R_0`$, to a second critical temperature $`T_20.04`$. Below $`T_2`$ the EI phase is stable with an almost constant $`\mathrm{T}_\mathrm{c}(|\mathrm{E}_\mathrm{g}|)`$, which, however, for larger band overlaps slowly decreases (see inset). The steeple-like shape of the phase diagram reflects the different phases from which the EI is approached: semi-conducting for $`\mathrm{E}_\mathrm{g}>0`$ and semi-metallic for $`\mathrm{E}_\mathrm{g}<0`$. Entering the EI phase from the semi-conductor side leads to formation of strongly bound excitons. On the other hand, when the EI phase is approached from the semi-metal, exciton formation is strongly suppressed due to the free carrier’s screening of the Coulomb potential. In that case, loosly bound Cooper-type pairs emerge, resulting in a rather fragile EI phase. The crossover from excitons to Cooper-type pairs occurs at $`\mathrm{E}_\mathrm{g}0.3`$, the band overlap, where the screening length becomes roughly equal to the exciton radius. For $`|\mathrm{E}_\mathrm{g}|4\mathrm{k}_\mathrm{B}\mathrm{T}`$, the critical temperature is exponentially small and approximately given by $`\mathrm{k}_\mathrm{B}\mathrm{T}_\mathrm{c}(\gamma |\mathrm{E}_\mathrm{g}|/\pi )\mathrm{exp}(\pi \sqrt{|\mathrm{E}_\mathrm{g}|}/\mathrm{ln}(1+\pi \sqrt{|\mathrm{E}_\mathrm{g}|}/2))`$, with $`\gamma =\mathrm{exp}(0.577)`$ (dashed line in the inset). Anisotropies in the band structure and other pair breaking effects would easily destroy this part of the phase diagram.
Support from SFB 652 is greatly acknowledged. |
warning/0506/hep-ph0506165.html | ar5iv | text | # Neutrino Physics1footnote 11footnote 1Lectures presented at the 2004 SLAC Summer Institute. To appear in the Proceedings. Fermilab publication FERMILAB-PUB-05-236-T.
## I PHYSICS OF NEUTRINO OSCILLATION
### I.1 Introduction
There has been a breakthrough in neutrino physics. It has been discovered that neutrinos have nonzero masses, and that leptons mix. The evidence for masses and mixing is the observation that neutrinos can change from one type, or “flavor”, to another. In this first section of these lectures, we will explain the physics of neutrino flavor change, or “oscillation”, as it is called. We will treat oscillation both in vacuum and in matter, and see why it implies masses and mixing.
That neutrinos have masses means that there is a spectrum of neutrino mass eigenstates $`\nu _i,i=1,2,\mathrm{}`$, each with a mass $`m_i`$. What leptonic mixing means may be understood by considering the leptonic decays $`W^+\nu _i+\overline{\mathrm{}_\alpha }`$ of the $`W`$ boson. Here, $`\alpha =e,\mu `$, or $`\tau `$, and $`\mathrm{}_e`$ is the electron, $`\mathrm{}_\mu `$ the muon, and $`\mathrm{}_\tau `$ the tau. The particle $`\mathrm{}_\alpha `$ is referred to as the charged lepton of flavor $`\alpha `$. Mixing means simply that in the $`W^+`$ decays to the particular charged lepton $`\overline{\mathrm{}_\alpha }`$, the accompanying neutrino mass eigenstate is not always the same $`\nu _i`$, but can be any of the different $`\nu _i`$. The amplitude for $`W^+`$ decay to produce the specific combination $`\overline{\mathrm{}_\alpha }+\nu _i`$ is denoted by $`U_{\alpha i}^{}`$. The neutrino state emitted in $`W^+`$ decay together with the particular charged lepton $`\overline{\mathrm{}_\alpha }`$ is then
$$|\nu _\alpha >=\underset{i}{}U_{\alpha i}^{}|\nu _i>.$$
(1)
This superposition of mass eigenstates is called the neutrino of flavor $`\alpha `$.
The quantities $`U_{\alpha i}`$ may be collected into a matrix known as the leptonic mixing matrix r1 . According to the Standard Model, $`U`$ is unitary. This unitarity guarantees that when the neutrino $`\nu _\alpha `$ interacts in a detector and creates a charged lepton, the latter will always be $`\mathrm{}_\alpha `$, the charged lepton with the same flavor as the neutrino. That is, a $`\nu _e`$ always yields an $`e`$, a $`\nu _\mu `$ a $`\mu `$, and a $`\nu _\tau `$ a $`\tau `$.
The relation (1), expressing a neutrino of definite flavor as a superposition of mass eigenstates, may be inverted to express each mass eigenstate $`\nu _i`$ as a superposition of flavors:
$$|\nu _i>=\underset{\alpha }{}U_{\alpha i}|\nu _\alpha >.$$
(2)
The flavor-$`\alpha `$ fraction of $`\nu _i`$ is obviously $`|U_{\alpha i}|^2`$. When $`\nu _i`$ interacts in a detector and produces a charged lepton, this flavor-$`\alpha `$ fraction is the probability that the charged lepton will be of flavor $`\alpha `$.
We turn now to the physics of neutrino oscillation.
### I.2 Neutrino Flavor Change in Vacuum
A typical neutrino flavor change, or “oscillation”, is depicted schematically in the top part of Fig. 1.
There, a neutrino source produces a neutrino together with a charged lepton $`\overline{\mathrm{}_\alpha }`$ of flavor $`\alpha `$. Thus, at birth, the neutrino is a $`\nu _\alpha `$. It then travels a distance $`L`$ to a detector. There, it interacts with a target and produces a second charged lepton $`\mathrm{}_\beta `$ of flavor $`\beta `$. Thus, at the time of its interaction in the detector, the neutrino is a $`\nu _\beta `$. If $`\beta \alpha `$ (for example, if $`\mathrm{}_\alpha `$ is a $`\mu `$ but $`\mathrm{}_\beta `$ is a $`\tau `$), then, during its journey to the detector, the neutrino has morphed from a $`\nu _\alpha `$ into a $`\nu _\beta `$.
This change of neutrino flavor, $`\nu _\alpha \nu _\beta `$, is a quintessentially quantum-mechanical effect. Indeed, it entails some quantum-mechanical subtleties that are still debated to this day r2 . However, there is little debate about the “ bottom line”—the expression for the flavor-change probability, P($`\nu _\alpha \nu _\beta `$). Therefore, in the interest of brevity, here we will derive this expression using an efficient approach r3 that contains all the essential quantum physics, even though it may not do justice to the subtleties.
Since, as described by Eq. (1), a $`\nu _\alpha `$ is actually a coherent superposition of mass eigenstates $`\nu _i`$, the particle that propagates from the neutrino source to the detector in Fig. 1 is one or another of the $`\nu _i`$, and we must add the contributions of the different $`\nu _i`$ coherently. Thus, with “Amp” denoting an amplitude, Amp($`\nu _\alpha \nu _\beta `$) is given by the lower part of Fig. 1. There, the contribution of each $`\nu _i`$ is a product of three factors. The first is the amplitude for the neutrino produced together with an $`\overline{\mathrm{}_\alpha }`$ by the source to be, in particular, a $`\nu _i`$. As has been said, this amplitude is $`U_{\alpha i}^{}`$. The second factor is the amplitude for the produced $`\nu _i`$ to propagate from the source to the detector. This factor is denoted by Prop($`\nu _i`$) in Fig. 1, and will be determined shortly. The final factor is the amplitude for the charged lepton created by the $`\nu _i`$ when it interacts in the detector to be, in particular, an $`\mathrm{}_\beta `$. From the Hermiticity of the Hamiltonian that describes neutrino-charged lepton-$`W`$ boson couplings, it follows that if Amp($`W\overline{\mathrm{}_\alpha }\nu _i)=U_{\alpha i}^{}`$, then Amp$`(\nu _i\mathrm{}_\beta W)=U_{\beta i}`$. Thus, the final factor in the $`\nu _i`$ contribution is $`U_{\beta i}`$, and
$$\mathrm{Amp}(\nu _\alpha \nu _\beta )=\underset{i}{}U_{\alpha i}^{}\mathrm{Prop}(\nu _i)U_{\beta i}.$$
(3)
Now, what is Prop($`\nu _i`$)? To find out, we go to the $`\nu _i`$ rest frame. We call the time in that frame $`\tau _i`$. If $`\nu _i`$ has rest mass $`m_i`$, then in its rest frame its state vector obeys the trivial Schrödinger equation
$$i\frac{}{\tau _i}|\nu _i(\tau _i)>=m_i|\nu _i(\tau _i)>.$$
(4)
The solution to this equation is obviously
$$|\nu _i(\tau _i)>=e^{im_i\tau _i}|\nu _i(0)>.$$
(5)
Thus, the amplitude for $`\nu _i`$ to propagate for a time $`\tau _i`$, which is just the amplitude $`<\nu _i(0)|\nu _i(\tau _i)>`$ for finding the original $`\nu _i`$ state $`|\nu _i(0)>`$ in the time evolved state $`|\nu _i(\tau _i)>`$, is simply $`\mathrm{exp}[im_i\tau _i]`$. Prop($`\nu _i`$) is just this amplitude with $`\tau _i`$ the proper time taken by $`\nu _i`$ to travel from the neutrino source to the detector.
For Prop($`\nu _i`$) to be useful to us, we must re-express it in terms of laboratory-frame variables. Two of these variables are the laboratory-frame distance, $`L`$, that the neutrino travels between its source and the detector, and the laboratory-frame time, $`t`$, that elapses during the trip. The value of $`L`$ is chosen by the experimenters through their choices for the location of the source and the location of the detector. Similarly, the value of $`t`$ is chosen by the experimenters through their choices for the time when the neutrino is created and the time when it is detected. Thus, $`L`$ and $`t`$ are defined by the experiment, and are common to all $`\nu _i`$ components of the beam. Different $`\nu _i`$ do not have different values of $`L`$ and $`t`$ from each other.
The other two laboratory-frame variables are the energy $`E_i`$ and momentum $`p_i`$ of mass eigenstate $`\nu _i`$ in the laboratory frame. By Lorentz invariance, the phase $`m_i\tau _i`$ in the $`\nu _i`$ propagator Prop($`\nu _i`$) is given in terms of the laboratory-frame variables by
$$m_i\tau _i=E_itp_iL.$$
(6)
The reader might object that, in reality, neutrino sources are essentially constant in time, and the time $`t`$ that elapses between the birth of a neutrino and its detection is not measured. This objection is quite correct. In practice, a realistic experiment averages over the time $`t`$ taken by the neutrino during its journey. Now, suppose two components of the neutrino beam, one with (laboratory-frame) energy $`E_1`$ and the other with energy $`E_2`$, contribute coherently to the neutrino signal observed in the detector. If the time taken by the neutrino to travel from source to detector is $`t`$, then by the time the beam component with energy $`E_j(j=1,2)`$ reaches the detector, it has picked up a phase factor $`\mathrm{exp}[iE_jt]`$. Thus, the interference between the $`E_1`$ and $`E_2`$ components of the beam will involve a phase factor $`\mathrm{exp}[i(E_1E_2)t]`$. Averaged over the unobserved travel time $`t`$, this factor vanishes, unless $`E_2=E_1`$. Thus, the only components of a neutrino beam that contribute coherently to a neutrino oscillation signal are components that have the same energy r4 . In particular, the different mass eigenstate components of a beam that contribute coherently to the oscillation signal must have the same energy, $`E`$.
At energy $`E`$, mass eigenstate $`\nu _i`$, with mass $`m_i`$, has a momentum $`p_i`$ given by
$$p_i=\sqrt{E^2m_i^2}E\frac{m_i^2}{2E}.$$
(7)
Here, we have used the fact that, given the extreme lightness of neutrinos, $`m_i^2E^2`$ for any realistic energy $`E`$. From Eqs. (6) and (7), we see that at energy $`E`$ the phase $`m_i\tau _i`$ in Prop($`\nu _i`$) is given by
$$m_i\tau _iE(tL)+\frac{m_i^2}{2E}L.$$
(8)
In this expression, the phase $`E(tL)`$ is irrelevant since it is common to all the interfering mass eigenstates. Thus, we may take
$$\mathrm{Prop}(\nu _i)=\mathrm{exp}[im_i^2\frac{L}{2E}].$$
(9)
Using this result, it follows from Eq. (3) that the amplitude for a neutrino to change from a $`\nu _\alpha `$ into a $`\nu _\beta `$ while traveling a distance $`L`$ through vacuum with energy $`E`$ is given by
$$\mathrm{Amp}(\nu _\alpha \nu _\beta )=\underset{i}{}U_{\alpha i}^{}e^{im_i^2\frac{L}{2E}}U_{\beta i}.$$
(10)
This expression holds for any number of flavors and mass eigenstates. Squaring it, we find that the probability P($`\nu _\alpha \nu _\beta `$) for $`\nu _\alpha \nu _\beta `$ is given by
$`\mathrm{P}(\nu _\alpha \nu _\beta )`$ $`=`$ $`|\mathrm{Amp}(\nu _\alpha \nu _\beta )|^2`$ (11)
$`=`$ $`\delta _{\alpha \beta }4{\displaystyle \underset{i>j}{}}\mathrm{}(U_{\alpha i}^{}U_{\beta i}U_{\alpha j}U_{\beta j}^{})\mathrm{sin}^2(\mathrm{\Delta }m_{ij}^2{\displaystyle \frac{L}{4E}})`$
$`+2{\displaystyle \underset{i>j}{}}\mathrm{}(U_{\alpha i}^{}U_{\beta i}U_{\alpha j}U_{\beta j}^{})\mathrm{sin}(\mathrm{\Delta }m_{ij}^2{\displaystyle \frac{L}{2E}}),`$
where
$$\mathrm{\Delta }m_{ij}^2m_i^2m_j^2.$$
(12)
In obtaining Eq. (11). we have made judicious use of the unitarity of $`U`$.
The oscillation probability P($`\nu _\alpha \nu _\beta `$) of Eq. (11) is that for a neutrino, rather than an antineutrino, as one can see from Fig. 1, which shows that the oscillating particle is produced in association with a charged antilepton $`\overline{\mathrm{}}`$, and produces a charged lepton $`\mathrm{}`$ in the detector. The probability P($`\overline{\nu _\alpha }\overline{\nu _\beta }`$) for the corresponding antineutrino oscillation may be found from P($`\nu _\alpha \nu _\beta `$) using the fact that the process $`\overline{\nu _\alpha }\overline{\nu _\beta }`$ is the CPT-mirror image of $`\nu _\beta \nu _\alpha `$. Thus, assuming that CPT invariance holds,
$$\mathrm{P}(\overline{\nu _\alpha }\overline{\nu _\beta })=\mathrm{P}(\nu _\beta \nu _\alpha ).$$
(13)
Now, from Eq. (11) we see that
$$\mathrm{P}(\nu _\beta \nu _\alpha ;U)=\mathrm{P}(\nu _\alpha \nu _\beta ;U^{}).$$
(14)
Thus, if CPT invariance holds, it follows from Eq. (11) that
P(
( )
[-.7ex] να
( )
[-.7ex] νβ
)P
( )
[-.7ex] να
( )
[-.7ex] νβ
\displaystyle\mathrm{P}(\shortstack{{\tiny(\rule[0.86108pt]{5.0pt}{0.28453pt})}
\\
[-.7ex] $\nu_{\alpha}$}\rightarrow\shortstack{{\tiny(\rule[0.86108pt]{5.0pt}{0.28453pt})}
\\
[-.7ex] $\nu_{\beta}$}) $`=`$ $`\delta _{\alpha \beta }4{\displaystyle \underset{i>j}{}}\mathrm{}(U_{\alpha i}^{}U_{\beta i}U_{\alpha j}U_{\beta j}^{})\mathrm{sin}^2(\mathrm{\Delta }m_{ij}^2{\displaystyle \frac{L}{4E}})`$ (15)
$`\text{+}\text{(}\text{}\text{)}\mathrm{\hspace{0.33em}2}{\displaystyle \underset{i>j}{}}\mathrm{}(U_{\alpha i}^{}U_{\beta i}U_{\alpha j}U_{\beta j}^{})\mathrm{sin}(\mathrm{\Delta }m_{ij}^2{\displaystyle \frac{L}{2E}}).`$
From these expressions, we see that if the mixing matrix $`U`$ is complex, P($`\overline{\nu _\alpha }\overline{\nu _\beta }`$) and P($`\nu _\alpha \nu _\beta `$) will in general differ. Since $`\overline{\nu _\alpha }\overline{\nu _\beta }`$ is the CP-mirror image of $`\nu _\alpha \nu _\beta `$, P$`(\overline{\nu _\alpha }\overline{\nu _\beta })\mathrm{P}(\nu _\alpha \nu _\beta )`$ would be a violation of CP invariance. So far, CP violation has been seen only in the quark sector, so its observation in neutrino oscillation would be most interesting.
With the expressions for P(
( )
[-.7ex] να
( )
[-.7ex] νβ
)
( )
[-.7ex] να
( )
[-.7ex] νβ
(\shortstack{{\tiny(\rule[0.86108pt]{5.0pt}{0.28453pt})}
\\
[-.7ex] $\nu_{\alpha}$}\rightarrow\shortstack{{\tiny(\rule[0.86108pt]{5.0pt}{0.28453pt})}
\\
[-.7ex] $\nu_{\beta}$}) in hand, let us note several features of neutrino oscillation:
1. If neutrinos are massless, so that all $`\mathrm{\Delta }m_{ij}^2=0`$, then, as we see from Eq. (15), P(
( )
[-.7ex] να
( )
[-.7ex] νβ
)=δαβ
( )
[-.7ex] να
( )
[-.7ex] νβ
subscript𝛿𝛼𝛽(\shortstack{{\tiny(\rule[0.86108pt]{5.0pt}{0.28453pt})}
\\
[-.7ex] $\nu_{\alpha}$}\rightarrow\shortstack{{\tiny(\rule[0.86108pt]{5.0pt}{0.28453pt})}
\\
[-.7ex] $\nu_{\beta}$})=\delta_{\alpha\beta}. Thus, the observation that neutrinos can change from one flavor to a different one implies neutrino mass. Indeed, it was this observation that led to the conclusion that neutrinos have nonzero masses.
To be sure, every flavor change seen so far has involved neutrinos that travel through matter. Equation (15) is for flavor change in vacuum, and does not take into account any interaction between the neutrinos and matter between their source and their detector. Thus, one might ask how we know that the observed flavor changes are not due to flavor-changing interactions between neutrinos and matter, rather than to neutrino masses. In response to this question, two points can be made. First, while the Standard Model of elementary particle physics does not include neutrino masses, it does include a rather well confirmed description of neutrino interactions, and this description states that neutrino interactions with matter do not change flavor. Secondly, for at least some of the observed flavor changes, matter effects are expected to be negligible, and there is evidence that for these cases, the flavor-change probability does depend on $`L`$ and $`E`$ in the combination $`L/E`$, as predicted by Eq. (15). Apart from a constant, $`L/E`$ is just the proper time that elapses in the rest frame of a neutrino as it travels a distance $`L`$ with energy $`E`$. Thus, these flavor changes appear to be an evolution of the neutrino itself over time, rather than a result of interaction with matter.
2. Suppose there is no leptonic mixing. This means that in the decay $`W^+\overline{\mathrm{}_\alpha }+\nu _i`$, which as we recall has an amplitude $`U_{\alpha i}^{}`$, the particular charged antilepton $`\overline{\mathrm{}_\alpha }`$ of flavor $`\alpha `$ is always accompanied by the same neutrino mass eigenstate $`\nu _i`$. That is, if $`U_{\alpha i}^{}0`$, then $`U_{\alpha j}`$ vanishes for all $`ji`$. Thus, as we see from Eq. (15), P(
( )
[-.7ex] να
( )
[-.7ex] νβ
)=δαβ
( )
[-.7ex] να
( )
[-.7ex] νβ
subscript𝛿𝛼𝛽(\shortstack{{\tiny(\rule[0.86108pt]{5.0pt}{0.28453pt})}
\\
[-.7ex] $\nu_{\alpha}$}\rightarrow\shortstack{{\tiny(\rule[0.86108pt]{5.0pt}{0.28453pt})}
\\
[-.7ex] $\nu_{\beta}$})=\delta_{\alpha\beta}. Hence, the observation that neutrinos can change from one flavor to a different one implies mixing.
3. One can detect neutrino flavor change in two ways. One way is to observe, in a beam of neutrinos which are initially all of flavor $`\alpha `$, the appearance of neutrinos of a new flavor $`\beta `$ that is different from the original flavor $`\alpha `$. This is called an appearance experiment. The other way is to start with a $`\nu _\alpha `$ beam of known flux, and to observe that some of this known $`\nu _\alpha `$ flux disappears. This is called a disappearance experiment.
4. Including the so-far-omitted factors of $`\mathrm{}`$ and $`c`$, the argument of the oscillatory quantity $`\mathrm{sin}^2[\mathrm{\Delta }m_{ij}^2L/4E]`$ that appears in Eq. (15) for P(
( )
[-.7ex] να
( )
[-.7ex] νβ
)
( )
[-.7ex] να
( )
[-.7ex] νβ
(\shortstack{{\tiny(\rule[0.86108pt]{5.0pt}{0.28453pt})}
\\
[-.7ex] $\nu_{\alpha}$}\rightarrow\shortstack{{\tiny(\rule[0.86108pt]{5.0pt}{0.28453pt})}
\\
[-.7ex] $\nu_{\beta}$}) is given by
$$\mathrm{\Delta }m_{ij}^2\frac{L}{4E}=1.27\mathrm{\Delta }m_{ij}^2(\mathrm{eV}^2)\frac{L(\mathrm{km})}{E(\mathrm{GeV})}.$$
(16)
Now, $`\mathrm{sin}^2[1.27\mathrm{\Delta }m_{ij}^2(\mathrm{eV}^2)L(\mathrm{km})/E(\mathrm{GeV})]`$ is appreciable so long as its argument is of order unity or larger. Thus, an experiment with a given $`L`$ (km) /$`E`$ (GeV) is sensitive to neutrino squared-mass splittings $`\mathrm{\Delta }m_{ij}^2(\mathrm{eV}^2)`$ all the way down to $`[L(\mathrm{km})/E(\mathrm{GeV}]^1`$. For example, an experiment with $`L10^4`$ km, the diameter of the earth, and $`E1`$ GeV is sensitive to $`\mathrm{\Delta }m_{ij}^2`$ down to $`10^4`$ eV<sup>2</sup>. As this illustrates, neutrino oscillation provides experimental access to very tiny neutrino masses. It does this by revealing quantum interferences between amplitudes whose relative phases are proportional to neutrino squared-mass differences, but can nevertheless be visibly large if $`L/E`$ is large enough.
5. As Eq. (15) shows, the probability of flavor change in vacuum oscillates with $`L/E`$. It is this behavior that has led neutrino flavor change to be called “neutrino oscillation”.
6. As Eq. (15) indicates, neutrino oscillation probabilities depend only on neutrino squared-mass splittings, and not on the individual squared neutrino masses themselves. Thus, as illustrated in Fig. 2, oscillation experiments can determine the neutrino squared-mass spectral pattern, but not how far above zero the entire pattern lies.
7. Neutrino flavor change does not alter the total flux in a neutrino beam, but merely redistributes it among the flavors. Indeed, from Eq. (15) and the unitarity of the $`U`$ matrix, it follows trivially that
βP(
( )
[-.7ex] να
( )
[-.7ex] νβ
)=1,subscript𝛽P
( )
[-.7ex] να
( )
[-.7ex] νβ
1\sum_{\beta}\mathrm{P}(\shortstack{{\tiny(\rule[0.86108pt]{5.0pt}{0.28453pt})}
\\
[-.7ex] $\nu_{\alpha}$}\rightarrow\shortstack{{\tiny(\rule[0.86108pt]{5.0pt}{0.28453pt})}
\\
[-.7ex] $\nu_{\beta}$})=1~{}~{}, (17)
where the sum is to encompass all flavors $`\beta `$, including the original flavor $`\alpha `$. Eq. (17) states that the probability that a neutrino changes its flavor, plus the probability that it does not do so, is unity. Hence, flavor change involves no change in total flux. However, some of the flavors $`\beta \alpha `$ into which a neutrino can transform itself might be sterile flavors; that is, flavors that do not enjoy normal weak interactions and consequently will not be detected in any feasible detector. If some of the original neutrino flux becomes sterile, then an experiment which measures the total active neutrino flux—that is, the sum of the $`\nu _e,\nu _\mu `$, and $`\nu _\tau `$ fluxes—will find it to be less than the original flux.
8. In the literature, treatments of neutrino oscillation frequently assume that the different mass eigenstates $`\nu _i`$ that contribute coherently to a beam have a common momentum, rather than the common energy that we have argued they must have. While the assumption of a common momentum is technically incorrect, it is a harmless error, since, as can easily be shown r7 , it leads to the same oscillation probabilities as we have found.
An important special case of the not-so-simple formula for P(
( )
[-.7ex] να
( )
[-.7ex] νβ
)
( )
[-.7ex] να
( )
[-.7ex] νβ
(\shortstack{{\tiny(\rule[0.86108pt]{5.0pt}{0.28453pt})}
\\
[-.7ex] $\nu_{\alpha}$}\rightarrow\shortstack{{\tiny(\rule[0.86108pt]{5.0pt}{0.28453pt})}
\\
[-.7ex] $\nu_{\beta}$}) in Eq. (15) is the case where only two different neutrinos are important. The two-neutrino approximation is a fairly accurate description of a number of experiments. Suppose, then, that only two mass eigenstates, which we shall call $`\nu _1`$ and $`\nu _2`$, and two corresponding flavor states, which we shall call $`\nu _e`$ and $`\nu _\mu `$, are significant. There is then only one squared-mass splitting, $`m_2^2m_1^2\mathrm{\Delta }m^2`$. Furthermore, omitting phase factors that can be shown to have no effect on oscillation, the mixing matrix $`U`$ takes the simple form
$$\begin{array}{ccc}& \nu _1\nu _2& \\ U=& \begin{array}{c}\nu _e\\ \nu _\mu \end{array}\left[\begin{array}{cc}\mathrm{cos}\theta & \mathrm{sin}\theta \\ \mathrm{sin}\theta & \mathrm{cos}\theta \end{array}\right].& \end{array}$$
(18)
Here, the symbols above and to the left of the matrix label its columns and rows. The $`U`$ of Eq. (18) is just a 2$`\times `$2 rotation matrix, and the rotation angle $`\theta `$ within it is referred to as the mixing angle. Inserting the $`U`$ of Eq. (18) and the single $`\mathrm{\Delta }m^2`$ into the general expression for P(
( )
[-.7ex] να
( )
[-.7ex] νβ
)
( )
[-.7ex] να
( )
[-.7ex] νβ
(\shortstack{{\tiny(\rule[0.86108pt]{5.0pt}{0.28453pt})}
\\
[-.7ex] $\nu_{\alpha}$}\rightarrow\shortstack{{\tiny(\rule[0.86108pt]{5.0pt}{0.28453pt})}
\\
[-.7ex] $\nu_{\beta}$}), Eq. (15), we immediately find that, for $`\beta \alpha `$, when only two neutrinos matter,
P(
( )
[-.7ex] να
( )
[-.7ex] νβ
)=sin22θsin2(Δm2L4E).P
( )
[-.7ex] να
( )
[-.7ex] νβ
superscript22𝜃superscript2Δsuperscript𝑚2𝐿4𝐸\mathrm{P}(\shortstack{{\tiny(\rule[0.86108pt]{5.0pt}{0.28453pt})}
\\
[-.7ex] $\nu_{\alpha}$}\rightarrow\shortstack{{\tiny(\rule[0.86108pt]{5.0pt}{0.28453pt})}
\\
[-.7ex] $\nu_{\beta}$})=\sin^{2}2\theta\sin^{2}(\Delta m^{2}\frac{L}{4E})~{}~{}. (19)
In addition, the probability that the neutrino does not change flavor is, as usual, unity minus the probability that it does change flavor.
### I.3 Neutrino Flavor Change in Matter
When an accelerator on the earth’s surface sends a beam of neutrinos several hundred kilometers to a waiting detector, the beam does not travel through a vacuum, but through earth matter. Coherent forward scattering of the neutrinos in the beam from particles they encounter along the way can then have a large effect. Assuming that the neutrino interactions with matter are the flavor-conserving ones described by the Standard Model, a neutrino in matter can undergo coherent forward scattering from ambient particles in two ways. First, if it is a $`\nu _e`$—and only if it is a $`\nu _e`$—it can exchange a $`W`$ boson with an electron. Coherent forward scattering by electrons via $`W`$ exchange gives rise to an extra interaction potential energy $`V_W`$ possessed by electron neutrinos in matter. Clearly, this extra energy from a lowest-order weak interaction will be proportional to $`G_F`$, the Fermi coupling constant. Equally clearly, this extra energy from $`\nu _ee`$ scattering will be proportional to $`N_e`$, the number of electrons per unit volume. From the Standard Model, we find that
$$V_W=+\sqrt{2}G_FN_e,$$
(20)
and that this interaction potential energy changes sign if we replace the $`\nu _e`$ in the beam by $`\overline{\nu _e}`$.
Secondly, a neutrino in matter can exchange a $`Z`$ boson with an ambient electron, proton, or neutron. The Standard Model tells us that any flavor of neutrino can do this, and that the amplitude for this $`Z`$ exchange is flavor independent. The Standard Model also tells us that, at zero momentum transfer, the $`Z`$ couplings to electron and proton are equal and opposite. Thus, assuming the matter through which our neutrino is traveling is electrically neutral (equal electron and proton densities), the electron and proton contributions to coherent forward neutrino scattering via $`Z`$ exchange will cancel out. Thus, the $`Z`$ exchange will give rise to a neutrino-flavor-independent extra interaction potential energy $`V_Z`$ that depends only on $`N_n`$, the number of neutrons per unit volume. From the Standard Model, we find that
$$V_Z=\frac{\sqrt{2}}{2}G_FN_n,$$
(21)
and that, as for $`V_W`$, this interaction potential energy changes sign if we replace the neutrinos in the beam by antineutrinos.
As already noted, Standard Model interactions do not change neutrino flavor. Thus, unless hypothetical non-Standard-Model flavor-changing interactions are at work, the observation of neutrino flavor change implies neutrino mass and mixing even when the neutrinos are passing through matter.
Neutrino propagation in matter is conveniently treated via a Schrödinger laboratory-frame time-evolution equation of the form
$$i\frac{}{t}|\nu (t)>=|\nu (t)>.$$
(22)
Here, $`|\nu (t)>`$ is a multi-component neutrino state vector, with one component for each of the possible neutrino flavors. Correspondingly, the Hamiltonian $``$ is a matrix in flavor space. To illustrate, let us describe the simple case where only two neutrino flavors, say $`\nu _e`$ and $`\nu _\mu `$, are important. Then
$$|\nu (t)>=\left[\begin{array}{c}f_e(t)\\ f_\mu (t)\end{array}\right],$$
(23)
where $`f_e(t)`$ is the amplitude for the neutrino to be a $`\nu _e`$ at time $`t`$, and similarly for $`f_\mu (t)`$. Correspondingly, $``$ is a 2$`\times `$2 matrix in $`\nu _e\nu _\mu `$ space.
It is instructive to work out the two-neutrino case first in vacuum, and then in matter. Using Eq. (1) for $`|\nu _\alpha >`$ in terms of the mass eigenstates, we find that the $`\nu _\alpha \nu _\beta `$ matrix element of the vacuum Hamiltonian, $`_{\mathrm{Vac}}`$, is given by
$`<\nu _\alpha |_{\mathrm{Vac}}|\nu _\beta >`$ $`=`$ $`<{\displaystyle \underset{i}{}}U_{\alpha i}^{}\nu _i|_{\mathrm{Vac}}|{\displaystyle \underset{j}{}}U_{\beta j}^{}\nu _j>`$ (24)
$`=`$ $`{\displaystyle \underset{j}{}}U_{\alpha j}U_{\beta j}^{}\sqrt{p^2+m_j^2}.`$
Here, we are assuming that our neutrino is in a beam of definite momentum $`p`$, common to all its mass eigenstate components. (As has already been remarked, this assumption, while technically incorrect, will lead us to the correct oscillation probability.) In the last step of Eq. (24), we have used the fact that $`_{\mathrm{Vac}}|\nu _j>=E_j|\nu _j>`$, where $`E_j=\sqrt{p^2+m_j^2}`$ is the energy of mass eigenstate $`\nu _j`$ at momentum $`p`$, and the fact that the mass eigenstates of the Hermitean Hamiltonian $`_{\mathrm{Vac}}`$ must be orthogonal.
As we have seen, neutrino flavor change is a quantum interference phenomenon. In this interference, only the relative phases of the interfering contributions matter. Thus, only the relative energies of these contributions, which will detemine their relative phases, matter. Consequently, if it is convenient, we may freely subtract from the Hamiltonian $``$ any multiple of the identity matrix $`I`$. This will not affect the differences between the eigenvalues of $``$, and thus it will not affect the predictions of $``$ for flavor change.
Of course, in the two-neutrino case there are only two mass eigenstates, $`\nu _1`$ and $`\nu _2`$, with one splitting $`\mathrm{\Delta }m^2m_2^2m_1^2`$ between them, and the $`U`$ matrix is given by Eq. (18). Inserting this matrix into Eq. (24), making the highly-relativistic approximation $`(p^2+m_j^2)^{1/2}p+m_j^2/2p`$, and subtracting from $`_{\mathrm{Vac}}`$ an irrelevant multiple of the identity matrix, we obtain
$$_{\mathrm{Vac}}=\frac{\mathrm{\Delta }m^2}{4E}\left[\begin{array}{cc}\mathrm{cos}2\theta & \mathrm{sin}2\theta \\ \mathrm{sin}2\theta & \mathrm{cos}2\theta \end{array}\right].$$
(25)
In writing this expression, we have used $`pE`$, where $`E`$ is the average energy of the neutrino mass eigenstates in our highly relativistic beam of momentum $`p`$.
It is easy to confirm that the Hamiltonian $`_{\mathrm{Vac}}`$ of Eq. (25) leads to the same two-neutrino oscillation probability, Eq. (19), as we have already found by other means. For example, consider the oscillation $`\nu _e\nu _\mu `$. From the first row of Eq. (18) for $`U`$,
$$|\nu _e>=|\nu _1>\mathrm{cos}\theta +|\nu _2>\mathrm{sin}\theta ,$$
(26)
while from the second row,
$$|\nu _\mu >=|\nu _1>\mathrm{sin}\theta +|\nu _2>\mathrm{cos}\theta .$$
(27)
Now, the eigenvalues of $`_{\mathrm{Vac}}`$, Eq. (25), are
$$\lambda _1=\frac{\mathrm{\Delta }m^2}{4E},\lambda _2=+\frac{\mathrm{\Delta }m^2}{4E}.$$
(28)
The corresponding eigenvectors, $`|\nu _1>`$ and $`|\nu _2>`$, are related to $`|\nu _e>`$ and $`|\nu _\mu >`$ by Eqs. (26) and (27). Thus, with $``$ the $`_{\mathrm{Vac}}`$ of Eq. (25), the Schrödinger equation of Eq. (22) implies that if at time $`t=0`$ we start with a $`|\nu _e>`$, then after a time $`t`$ this $`|\nu _e>`$ will evolve into the state
$$|\nu (t)>=|\nu _1>e^{+i\frac{\mathrm{\Delta }m^2}{4E}t}\mathrm{cos}\theta +|\nu _2>e^{i\frac{\mathrm{\Delta }m^2}{4E}t}\mathrm{sin}\theta .$$
(29)
The probability P$`(\nu _e\nu _\mu )`$ that this time-evolved neutrino will be detected as a $`\nu _\mu `$ is then, from Eqs. (27) and (29),
$`\mathrm{P}(\nu _e\nu _\mu )`$ $`=`$ $`|<\nu _\mu |\nu (t)>|^2`$ (30)
$`=`$ $`|\mathrm{sin}\theta \mathrm{cos}\theta (e^{i\frac{\mathrm{\Delta }m^2}{4E}t}+e^{i\frac{\mathrm{\Delta }m^2}{4E}t})|^2`$
$`=`$ $`\mathrm{sin}^22\theta \mathrm{sin}^2(\mathrm{\Delta }m^2{\displaystyle \frac{L}{4E}}).`$
In the last step, we have replaced our highly-relativistic neutrino’s travel time $`t`$ by its travel distance $`L`$. The flavor change probability of Eq. (30) does indeed agree with what we found earlier, Eq. (19).
Let us turn now to neutrino propagation in matter. There, the 2$`\times `$2 vacuum Hamiltonian $`_{\mathrm{Vac}}`$ is replaced by a matrix $`_M`$ given by
$$_M=_{\mathrm{Vac}}+V_W\left[\begin{array}{cc}1& 0\\ 0& 0\end{array}\right]+V_Z\left[\begin{array}{cc}1& 0\\ 0& 1\end{array}\right].$$
(31)
Here, the second term on the right-hand side is the contribution from the interaction potential energy caused by $`W`$ exchange, Eq. (20). Since this energy affects only $`\nu _e`$, its contribution is nonvanishing only in the upper left, $`\nu _e\nu _e`$, element of $`_M`$. The last term on the right-hand side of Eq. (31) is the contribution from the interaction potential energy caused by $`Z`$ exchange, Eq. (21). Since this energy affects all flavors equally, its contribution to $`_M`$ is a multiple of the identity matrix, and consequently can be dropped. Then
$$_M=_{\mathrm{Vac}}+\frac{V_W}{2}\left[\begin{array}{cc}1& 0\\ 0& 1\end{array}\right]+\frac{V_W}{2}\left[\begin{array}{cc}1& 0\\ 0& 1\end{array}\right],$$
(32)
where we have now split the $`W`$-exchange contribution into a piece that is not proportional to the identity, plus a piece that is proportional to it. Dropping the irrelevant latter piece as well, we have from Eqs. (25) and (32)
$$_M=\frac{\mathrm{\Delta }m^2}{4E}\left[\begin{array}{cc}(\mathrm{cos}2\theta x)& \mathrm{sin}2\theta \\ \mathrm{sin}2\theta & (\mathrm{cos}2\theta x)\end{array}\right],$$
(33)
in which
$$x\frac{V_W/2}{\mathrm{\Delta }m^2/4E}=\frac{2\sqrt{2}G_FN_eE}{\mathrm{\Delta }m^2}.$$
(34)
The parameter $`x`$ is a measure of the importance of the matter effect relative to that of the neutrino squared-mass splitting.
If we define
$$\mathrm{\Delta }m_M^2\mathrm{\Delta }m^2\sqrt{\mathrm{sin}^22\theta +(\mathrm{cos}2\theta x)^2}$$
(35)
and
$$\mathrm{sin}^22\theta _M\frac{\mathrm{sin}^22\theta }{\mathrm{sin}^22\theta +(\mathrm{cos}2\theta x)^2},$$
(36)
then $`_M`$ can be written as
$$_M=\frac{\mathrm{\Delta }m_M^2}{4E}\left[\begin{array}{cc}\mathrm{cos}2\theta _M& \mathrm{sin}2\theta _M\\ \mathrm{sin}2\theta _M& \mathrm{cos}2\theta _M\end{array}\right].$$
(37)
That is, the Hamiltonian in matter, $`_M`$, is identical to its vacuum counterpart, $`_{\mathrm{Vac}}`$, Eq. (25), except that the vacuum parameters $`\mathrm{\Delta }m^2`$ and $`\theta `$ are replaced, respectively, by $`\mathrm{\Delta }m_M^2`$ and $`\theta _M`$.
Needless to say, the eigenstates of $`_M`$ differ from their vacuum counterparts. The splitting between the effective squared-masses of these eigenstates in matter differs from the vacuum splitting $`\mathrm{\Delta }m^2`$, and the effective mixing angle in matter—the angle that determines the $`\nu _e,\nu _\mu `$ composition of the eigenstates in matter—differs from the vacuum mixing angle $`\theta `$. Now, all of the physics of neutrino propagation in matter is contained in the matter Hamiltonian $`_M`$. But, according to Eq. (37), $`_M`$ depends on the parameters $`\mathrm{\Delta }m_M^2`$ and $`\theta _M`$ in exactly the same way as the vacuum Hamiltonian $`_{\mathrm{Vac}}`$, Eq. (25), depends on $`\mathrm{\Delta }m^2`$ and $`\theta `$. Thus, $`\mathrm{\Delta }m_M^2`$ must be the splitting between the effective squared-masses of the eigenstates in matter, and $`\theta _M`$ must be the effective mixing angle in matter.
In an experiment where an accelerator-generated neutrino beam is sent to a detector that is, say, 1000 km away, the beam passes through earth matter, but does not penetrate very deeply into the earth. The matter density encountered by such a beam en route is very roughly constant. Thus, the electron density $`N_e`$, hence the parameter $`x`$, hence the matter Hamiltonian $`_M`$, is roughly position independent, just like the vacuum Hamiltonian $`_{\mathrm{Vac}}`$. Comparing Eqs. (37) and (25), we then see that since $`_{\mathrm{Vac}}`$ leads to the vacuum oscillation probability P$`(\nu _e\nu _\mu )`$ of Eq. (30), $`_M`$ must lead to the in-matter oscillation probability
$$\mathrm{P}_M(\nu _e\nu _\mu )=\mathrm{sin}^22\theta _M\mathrm{sin}^2(\mathrm{\Delta }m_M^2\frac{L}{4E}).$$
(38)
That is, the oscillation probability in matter is the same as in vacuum, except for the replacement of the vacuum parameters $`\theta `$ and $`\mathrm{\Delta }m^2`$ by their in-matter equivalents.
How large is the earth matter effect, and what are its consequences? To answer this question, we first note from Eq. (34) that the parameter $`x`$, which measures the relative importance of matter, is proportional to the neutrino energy $`E`$. To estimate the proportionality constant, let us imagine that we have an accelerator-generated neutrino beam that travels $``$1000 km between its source and its detector. The electron density $`N_e`$ encountered by such a beam will be that of the earth’s mantle. The splitting $`\mathrm{\Delta }m^2`$ that will dominate the behavior of such a beam will be the “atmospheric” $`\mathrm{\Delta }m^2`$ that also governs the behavior of atmospheric neutrinos, and whose size is approximately $`2.4\times 10^3`$eV<sup>2</sup> r8 . Then from Eq. (34)
$$|x|\frac{E}{12\mathrm{GeV}}.$$
(39)
Thus, in a beam with $`E`$, say, 2 GeV, the matter effect is modest but not negligible, while in a beam with $`E`$, say, 20 GeV, the matter effect is very large.
We recall that the splitting $`\mathrm{\Delta }m^2`$ which appears in Eq. (34) is defined as $`m_2^2m_1^2`$, so that, depending on whether $`\nu _2`$ is heavier or lighter than $`\nu _1`$, $`\mathrm{\Delta }m^2`$ is positive or negative. We also recall that if the neutrinos, whose propagation in matter we have treated explicitly, are replaced by antineutrinos, then the interaction potential energy $`V_W`$, which is positive for neutrinos, reverses sign. As a result of these two effects, the sign of $`x`$, which for neutrinos is given by Eq. (34), is as summarized in Table 1.
Now, Eqs. (35) and (36) hold for both neutrinos and antineutrinos, so long as the appropriate value and sign of $`x`$ are used. These equations show that the effective squared-mass splitting in matter, $`\mathrm{\Delta }m_M^2`$, and the effective mixing angle in matter, $`\theta _M`$, both depend on the sign of $`x`$. Thus, since $`x`$(antineutrinos) = -$`x`$(neutrinos), $`\mathrm{\Delta }m_M^2`$ and $`\theta _M`$ will have different values for antineutrinos than they do for neutrinos. That is, there will be an asymmetry between antineutrino oscillation and neutrino oscillation that is induced by matter effects. This asymmetry has nothing to do with genuine CP violation, and will have to be disentangled from the antineutrino-neutrino asymmetry that does come from genuine CP violation in order for us to be able to study the latter phenomenon. However, the antineutrino-neutrino asymmetry coming from matter effects is by no means all bad. If the sign of $`\mathrm{cos}2\theta `$ is known, then we can use Eqs. (35) and (36), applied to both neutrinos and antineutrinos, to learn whether $`x`$(neutrinos) is positive and $`x`$(antineutrinos) is negative, or the other way around. This, in turn, will tell us whether the neutrino we have called $`\nu _2`$ is heavier or lighter than the one we have called $`\nu _1`$. As we shall see in Sec. II.3, it is hoped that a technique of precisely this kind can tell us whether the three-neutrino spectrum that may actually describe nature has the character shown in Fig. 2, or an inverted character, with the closely-spaced pair at the top, rather than at the bottom.
In principle, matter effects can have very dramatic consequences. From Eq. (36) for the effective mixing angle in matter, $`\theta _M`$, we see that if the vacuum mixing angle $`\theta `$ is tiny, with, say, $`\mathrm{sin}^22\theta =10^3`$, but $`x\mathrm{cos}2\theta `$, then $`\mathrm{sin}^22\theta _M`$ can be near or at its maximum possible value, unity. This dramatic amplification of a tiny mixing angle in vacuum into a very large one in matter is the “resonant” version of the Mikheyev-Smirnov-Wolfenstein effect r9 . It used to be thought that this dramatic amplification is actually occurring inside the sun. However, we now know that the solar neutrino mixing angle is already quite large ($`34^{}`$) in vacuum r10 . Thus, while the effect of solar matter on solar neutrinos is still very significant, it is not quite as dramatic as once thought.
## II WHAT WE HAVE LEARNED AND THE OPEN QUESTIONS
Now that we have discussed the physics of neutrino oscillation, let us consider what the experimental evidence for oscillation has taught us about neutrinos, and the questions that this evidence has raised.
### II.1 What Have We Learned?
We do not yet know how many neutrino mass eigenstates there are. Of course, there are at least three, but if the Liquid Scintillator Neutrino Detector (LSND) experiment is confirmed, then there must be more than three, or else our usual assumptions about the neutrinos are wrong in some even more extreme way. On the other hand, if LSND is not confirmed, then Nature may contain only three neutrinos. In that case, the neutrino oscillation data (excluding LSND) tell us that the neutrino squared-mass spectrum is as shown in Fig. 3. There are a pair of mass eigenstates, $`\nu _1`$ and $`\nu _2`$, separated by the splitting $`\mathrm{\Delta }m_{21}^2m_2^2m_1^2\mathrm{\Delta }m_{\mathrm{sol}}^28.0\times 10^5`$eV<sup>2</sup>. This splitting drives the behavior of solar neutrinos. In addition, there is an isolated mass eigenstate, $`\nu _3`$, separated from $`\nu _2`$ and $`\nu _1`$ (“the solar pair”) by the splitting $`\mathrm{\Delta }m_{32}^2m_3^2m_2^3\mathrm{\Delta }m_{\mathrm{atm}}^22.4\times 10^3`$eV<sup>2</sup>. This splitting drives the behavior of atmospheric neutrinos.
As explained in Sec. I.1, each mass eigenstate is a superposition of flavors, and Fig. 3 shows its approximate flavor content \[see the figure caption\].
As mentioned in Sec. I.1, when mass eigenstate $`\nu _i`$ interacts and produces a charged lepton, the probability that this charged lepton is, in particular, an electron, is the $`\nu _e`$ fraction of $`\nu _i`$, $`|U_{ei}|^2`$. The probability that the charged lepton is a muon is the $`\nu _\mu `$ fraction of $`\nu _i`$, $`|U_{\mu i}|^2`$, and the probability that it is a tau is the $`\nu _\tau `$ fraction, $`|U_{\tau i}|^2`$.
Fig. 3 summarizes how we learned the flavor content of the various mass eigenstates, and the squared-mass splittings between them. With reference to this figure, let us explain how these features of the neutrino spectrum were found, starting with $`\nu _3`$.
The $`\nu _e`$ fraction of $`\nu _3`$ is not known, but is bounded by reactor experiments that had a detector at a distance $`L1`$ km from the reactor. Since the (anti)neutrinos emitted by a reactor have an energy $`E3`$ GeV, this detector distance made these experiments sensitive to oscillation involving the larger (mass)<sup>2</sup> gap, $`\mathrm{\Delta }m_{\mathrm{atm}}^22.4\times 10^3`$ eV<sup>2</sup>, but not to oscillation involving the smaller gap, $`\mathrm{\Delta }m_{\mathrm{sol}}^28.0\times 10^5`$ eV<sup>2</sup> \[cf. Eq. (16) and surrounding text\]. As a result, these experiments probed the properties of $`\nu _3`$, the isolated neutrino at one end of the $`\mathrm{\Delta }m_{\mathrm{atm}}^2`$ gap r11 . In particular, they probed the $`\nu _e`$ fraction of $`\nu _3`$, since the particles emitted by a reactor are $`\overline{\nu _e}`$. The experiments saw no oscillation of these $`\overline{\nu _e}`$, whose disappearance they sought, and thereby set a 3$`\sigma `$ upper bound of $`|U_{e3}|^2<0.045`$ on the $`\nu _e`$ fraction of $`\nu _3`$ r12 .
One hears a lot of discussion of a leptonic mixing angle called $`\theta _{13}`$. This angle is so defined that $`|U_{e3}|^2=\mathrm{sin}^2\theta _{13}`$. Thus, $`\theta _{13}`$ is a measure of the smallness of the $`\nu _e`$ part of $`\nu _3`$.
Apart from this small $`\nu _e`$ piece, $`\nu _3`$ is of $`\nu _\mu `$ and $`\nu _\tau `$ flavor. Now, the oscillation of atmospheric muon neutrinos is observed to be dominated by $`\nu _\mu \nu _\tau `$, with a $`\nu _\mu \nu _\tau `$ mixing angle that is very large. The best fit for this angle is maximal mixing: 45. This atmospheric mixing angle will be reflected in the flavor content of $`\nu _3`$, since $`\nu _3`$ is at one end of the splitting $`\mathrm{\Delta }m_{\mathrm{atm}}^2`$ that drives atmospheric neutrino oscillation. If the angle is truly maximal, then, apart from its small $`\nu _e`$ component, $`\nu _3`$ is simply $`(\nu _\mu +\nu _\tau )/\sqrt{2}`$. This mimics the behavior of the neutral $`K`$ meson system. There, apart from a small CP violation, the mixing of $`K^0`$ and $`\overline{K^0}`$ is maximal, with the consequence that $`K_S=(K^0+\overline{K^0})/\sqrt{2}`$.
It is found that 1-GeV upward-going atmospheric neutrinos, which originate in the atmosphere on the far side of the earth from their detector, and hence travel $`10^4`$ km—the diameter of the earth—to reach the detector, oscillate. In contrast, 1-GeV downward-going neutrinos, which originate in the atmosphere directly above their detector, and thus travel only $`10`$ km—the thickness of the atmosphere—to reach the detector, do not oscillate. From these facts, and Eq. (16), it follows that $`10^4`$eV$`{}_{}{}^{2}\stackrel{<}{}\mathrm{\Delta }m_{\mathrm{atm}}^2\stackrel{<}{}\mathrm{\hspace{0.33em}10}^2`$eV<sup>2</sup>. Of course, the detailed experiments pin down $`\mathrm{\Delta }m_{\mathrm{atm}}^2`$ much better than that.
It is well established that the change of flavor of solar neutrinos is caused by the Large Mixing Angle version of the Mikheyev-Smirnov-Wolfenstein effect (LMA-MSW). As will be explained shortly, according to LMA-MSW, the probability P$`{}_{\mathrm{sol}}{}^{}(\nu _e\nu _e)`$ that solar neutrinos, all of which are born $`\nu _e`$, will still be $`\nu _e`$ when they arrive at the earth, is to a good approximation just the $`\nu _e`$ fraction of the mass eigenstate $`\nu _2`$. Experimentally, P$`{}_{\mathrm{sol}}{}^{}(\nu _e\nu _e)`$ is found to be approximately 1/3 r10 . Therefore, the $`\nu _e`$ fraction of $`\nu _2`$, $`|U_{e2}|^2`$, is $``$ 1/3. From the unitarity of the $`U`$ matrix, $`|U_{e1}|^2+|U_{e2}|^2+|U_{e3}|^2=1`$. Thus, since $`|U_{e3}|^2`$ is small and $`|U_{e2}|^2`$ 1/3, $`|U_{e1}|^2`$, the $`\nu _e`$ fraction of $`\nu _1`$, is approximately 2/3.
If the atmospheric $`\nu _\mu \nu _\tau `$ mixing is maximal, then, together, $`\nu _1`$ and $`\nu _2`$ contain the neutrino $`(\nu _\mu \nu _\tau )/\sqrt{2}`$, the 50-50 $`\nu _\mu \nu _\tau `$ mixture that is orthogonal to the one occurring in $`\nu _3`$. Then, $`\nu _1,\nu _2`$, and $`\nu _3`$ all contain $`\nu _\mu `$ and $`\nu _\tau `$ in equal proportion.
The solar (mass)<sup>2</sup> splitting, $`\mathrm{\Delta }m_{\mathrm{sol}}^2`$, has been determined from the observed energy dependence of solar neutrino flavor change, and especially from the observed energy dependence of reactor $`\overline{\nu _e}`$ disappearance as studied by the KamLAND reactor experiment r14 . Unlike the reactor experiments alluded to earlier, KamLAND has a detector situated, not 1 km, but on average 180 km away from reactors that are its $`\overline{\nu _e}`$ sources. With this greatly increased source-to-detector distance, KamLAND is sensitive to oscillation involving the small splitting $`\mathrm{\Delta }m_{\mathrm{sol}}^2`$.
While the flavor content pictured in Fig. 3 tells us the magnitudes $`|U_{\alpha i}|^2`$, it does not show the signs or phases of the $`U`$ matrix elements. To discuss these, we need to look directly at $`U`$ itself. It is illuminating to write $`U`$ in the form
$`\mathrm{Atmospheric}\text{Cross-Mixing}\mathrm{Solar}\text{Majorana CP-violating phases}`$
$`U`$ $`=`$ $`\left[\begin{array}{ccc}1& 0& 0\\ 0& c_{23}& s_{23}\\ 0& s_{23}& c_{23}\end{array}\right]\times \left[\begin{array}{ccc}c_{13}& 0& s_{13}e^{i\delta }\\ 0& 1& 0\\ s_{13}e^{i\delta }& 0& c_{13}\end{array}\right]\times \left[\begin{array}{ccc}c_{12}& s_{12}& 0\\ s_{12}& c_{12}& 0\\ 0& 0& 1\end{array}\right]\times \left[\begin{array}{ccc}e^{i\alpha _1/2}& 0& 0\\ 0& e^{i\alpha _2/2}& 0\\ 0& 0& 1\end{array}\right].`$ (52)
Here, $`c_{ij}\mathrm{cos}\theta _{ij}`$ and $`s_{ij}\mathrm{sin}\theta _{ij}`$, where $`\theta _{ij}`$ is a mixing angle, and $`\delta ,\alpha _1`$, and $`\alpha _2`$ are CP-violating phases.
In the decomposition of $`U`$ in Eq. (52), the first, Atmospheric, factor matrix dominates the mixing exhibited by the atmospheric neutrinos. This factor depends on a mixing angle $`\theta _{23}`$ that is approximately the atmospheric mixing angle $`\theta _{\mathrm{atm}}`$ determined when the atmospheric $`\nu _\mu \nu _\tau `$ oscillation is described by an approximate two-neutrino oscillation formula of the form of Eq. (19). At 90% confidence level, $`37^{}\theta _{\mathrm{atm}}53^{}`$ r8 . As already mentioned, this mixing is very large. The value of $`\theta _{\mathrm{atm}}`$ that fits the data best is 45: maximal mixing r8 .
The third, Solar, matrix on the right-hand side of Eq. (52) dominates the mixing of solar neutrinos. The mixing angle in this matrix, $`\theta _{12}`$, is approximately the solar mixing angle $`\theta _{\mathrm{sol}}`$ determined by approximating the solar neutrino problem as a two-neutrino problem. This is a very good approximation, since the solar neutrinos are born as $`\nu _e`$, and $`|U_{e3}|^2`$, the $`\nu _3`$ fraction of $`\nu _e`$, is quite small. Thus, $`\nu _e`$ is approximately a mixture of just $`\nu _1`$ and $`\nu _2`$, and it will remain so as it propagates. Experimentally, $`\theta _{\mathrm{sol}}=(33.9\genfrac{}{}{0pt}{}{+2.4}{2.2})`$ degrees r10 . Although this mixing angle is not maximal, as $`\theta _{\mathrm{atm}}`$ might be, it is still very large. It is interesting that two of the leptonic mixing angles are large, while all of the quark mixing angles are small. Is there a clue to the physics underlying fermion mass and mixing in this disparity?
The second, Cross-Mixing, matrix in Eq. (52) involves the small mixing angle $`\theta _{13}`$, which measures the small $`\nu _e`$ part of $`\nu _3`$. This matrix also may involve a CP-violating phase $`\delta `$. A nonvanishing $`\delta `$ will in general lead to a CP-volating difference between the probabilities for the CP-mirror-image oscillations $`\nu _\alpha \nu _\beta `$ and $`\overline{\nu _\alpha }\overline{\nu _\beta }`$. We note from Eq. (52), however, that $`\delta `$ enters the mixing matrix $`U`$ only in combination with $`\mathrm{sin}\theta _{13}`$. Thus, the size of the CP-violating difference $`P(\nu _\alpha \nu _\beta )P(\overline{\nu _\alpha }\overline{\nu _\beta })`$ will depend on $`\theta _{13}`$. Consequently, the power of the facilities we will need in order to see this CP-violating difference will also depend, at least roughly, on $`\theta _{13}`$. At present, we know only that, at 3$`\sigma `$, $`\theta _{13}<12^{}`$ r12 . Clearly, demonstrating that $`\theta _{13}`$ is nonvanishing, and determining its approximate size, are rather important.
The final matrix in Eq. (52) contains so-called “Majorana” CP-violating phases that have no analogue in the quark sector, and that have physical effects only if the neutrinos are their own antiparticles. In particular, these phases have no effect on neutrino oscillation, which is completely insensitive to whether neutrinos are their own antiparticles.
The reader may wonder why the CP-violating phase $`\delta `$ enters $`U`$ in combination with $`\theta _{13}`$, rather than $`\theta _{12}`$ or $`\theta _{23}`$. Doesn’t the $`U`$ matrix decomposition of Eq. (52) treat all three mixing angles symmetrically? The answer to this puzzle is that if one multiplies the matrices in Eq. (52) to obtain the detailed expression for $`U`$, and then uses one’s freedom to rephase elements of $`U`$ by phase-redefining any neutrino or charged lepton, then one can completely remove the phase $`\delta `$ from $`U`$ if any of the three mixing angles vanishes. That is, for $`\delta `$ to have physical consequences, all three of the mixing angles must be non-zero. However, we already know that two of them, $`\theta _{12}`$ and $`\theta _{23}`$, are very far from zero. Therefore, we have chosen to write $`U`$ in a form that emphasizes the fact that the CP-violating consequences of $`\delta `$ will disappear if the remaining mixing angle $`\theta _{13}`$—the only mixing angle we don’t know—should vanish.
Figures 5 and 5 show graphically the allowed regions for the atmospheric and solar mixing angles and (mass)<sup>2</sup> splittings.
### II.2 The Large Mixing Angle MSW Effect
In Sec. I.3, we discussed neutrino flavor change in matter. Let us now consider how this works in the case of the solar neutrinos, which encounter a lot of matter as they travel from the core of the sun, where they are born, to its outer edge. As mentioned previously, it is now well established that the change of flavor of these neutrinos is governed by the Large Mixing Angle MSW Effect. As we saw from Eq. (34), the importance of matter effects grows with neutrino energy. In the sun, this has the consequence that for the neutrinos produced in <sup>8</sup>B decay, which are at the high end of the solar neutrino energy spectrum, the influence of solar matter is quite significant. However, for the neutrinos produced by $`pp`$ fusion, which have an energy an order of magnitude lower, the effect of matter is tiny. Since it is the <sup>8</sup>B neutrinos that have been most extensively studied, let us focus on them.
As we have already noted, solar neutrino propagation is approximately a two-neutrino phenomenon. Solar neutrinos are mixtures of $`\nu _1`$ and $`\nu _2`$, with a $`\nu _3`$ component that may be neglected. The two mass eigenstates $`\nu _1`$ and $`\nu _2`$ make up two flavor states: $`\nu _e`$ and a state $`\nu _x`$ that is a coherent linear combination of $`\nu _\mu `$ and $`\nu _\tau `$. The precise $`\nu _\mu \nu _\tau `$ composition of $`\nu _x`$ depends on the $`\nu _\mu \nu _\tau `$ mixing angle $`\theta _{\mathrm{atm}}`$ measured in atmospheric neutrino oscillation. For $`\theta _{\mathrm{atm}}=45^{},\nu _x=(\nu _\mu \nu _\tau )/\sqrt{2}`$, a 50-50 mixture. Solar neutrino flavor change is the process $`\nu _e\nu _x`$.
Since solar neutrino propagation involves only two neutrinos, we may describe it by the two-neutrino in-matter Hamiltonian $`_M`$ of Eq. (31). Dropping the irrelevant $`Z`$ exchange term as before, and using Eqs. (25) and (20), we have
$$_M=\frac{\mathrm{\Delta }m_{\mathrm{sol}}^2}{4E}\left[\begin{array}{cc}\hfill \mathrm{cos}2\theta _{\mathrm{sol}}& \hfill \mathrm{sin}2\theta _{\mathrm{sol}}\\ \hfill \mathrm{sin}2\theta _{\mathrm{sol}}& \hfill \mathrm{cos}2\theta _{\mathrm{sol}}\end{array}\right]+\sqrt{2}G_FN_e\left[\begin{array}{cc}1& 0\\ 0& 0\end{array}\right].$$
(53)
The matrices in this equation are in the $`\nu _e\nu _x`$ flavor basis, and we have specialized to the case of solar neutrinos, where the relevant vacuum squared-mass splitting is $`\mathrm{\Delta }m_{\mathrm{sol}}^2`$ and the relevant vacuum mixing angle is $`\theta _{\mathrm{sol}}`$.
At the center of the sun, where the solar neutrinos are created, the coefficient of the second, matter-interaction term in the $`_M`$ of Eq. (53) is
$$\sqrt{2}G_FN_e0.75\times 10^5\mathrm{eV}^2/\mathrm{MeV}.$$
(54)
By comparison, since $`\mathrm{\Delta }m_{\mathrm{sol}}^28\times 10^5`$eV<sup>2</sup>, at a typical <sup>8</sup>B neutrino energy $`E`$ of $``$8 MeV, the coefficient of the first, vacuum term in this $`_M`$ is
$$\frac{\mathrm{\Delta }m_{\mathrm{sol}}^2}{4E}0.25\times 10^5\mathrm{eV}^2/\mathrm{MeV}.$$
(55)
Thus, at $`r=0`$ (where $`r`$ is the distance from the center of the sun), the interaction term in $`_M`$ dominates, at least to some extent. Let us then, as a first approximation, neglect the vacuum term in $`_M(r=0)`$. As we see from Eq. (53), $`_M(r=0)`$ is then diagonal. This means that a $`\nu _e`$ born at $`r=0`$ is an eigenstate of $`_M(r=0)`$. Since, from Eq. (53), the $`\nu _e`$ eigenvalue of (the approximate) $`_M(r=0)`$ is $`\sqrt{2}G_FN_e`$, while the other eigenvalue is zero, our $`\nu _e`$ is born in the higher-energy eigenstate of $`_M(r=0)`$.
Under the conditions where the Large Mixing Angle MSW effect occurs, the propagation of a <sup>8</sup>B neutrino from $`r0`$ to the outer edge of the sun is adiabatic. That is, the electron density $`N_e`$ changes slowly enough with $`r`$ that the Schrödinger equation, Eq. (22), in which $`_M`$ acts may be solved for one radius $`r`$ at at time, and then the solutions may be patched together. Our neutrino will propagate outward as the slowly-changing eigenstate of the slowly changing $`_M(r)`$. Furthermore, it can easily be shown that, as $`N_e`$ changes with $`r`$, the eigenvalues of $`_M(r)`$, Eq. (53), never cross. Thus, our neutrino, having begun as the higher-energy eigenstate of $`_M(r=0)`$, will emerge from the outer edge of the sun as the higher-energy eigenstate of $`_M(r=`$ radius of sun). But $`_M(r=`$ radius of sun) is just the vacuum Hamiltonian $`_{\mathrm{Vac}}`$, since at the outer edge of the sun, the electron density $`N_e`$ has fallen to zero. Thus, our neutrino emerges from the sun as the heavier mass eigenstate of the two-neutrino vacuum Hamiltonian. This is the mass eigenstate we have called $`\nu _2`$ in Fig. 3.
Being a mass eigenstate, our $`\nu _2`$ is an ordinary particle, just like an electron, and it will propagate from the outer edge of the sun to the surface of the earth without changing into anything else. The probability that when it interacts here on earth in a detector, it will do so like a $`\nu _e`$, producing an electron, is then just the $`\nu _e`$ fraction of $`\nu _2`$, $`|U_{e2}|^2`$. The Sudbury Neutrino Observatory (SNO) has found this probability to be approximately one-third r10 . Therefore, $`|U_{e2}|^21/3`$. (In the approximation where we treat the solar neutrinos neglecting $`\nu _3`$, the $`U`$ matrix is given by Eq. (18), with $`\theta `$ the solar mixing angle $`\theta _{\mathrm{sol}}`$, and with $`\nu _\mu `$ replaced by $`\nu _x`$. Then $`|U_{e2}|^2=\mathrm{sin}^2\theta _{\mathrm{sol}}`$, so the SNO results imply that $`\mathrm{sin}^2\theta _{\mathrm{sol}}1/3`$.)
This is how the Large Mixing Angle MSW effect works, at least to a good approximation. It should be noted that, while this effect causes a change in the flavor content of a neutrino, there is no sinusoidal oscillation of the sort that characterizes the flavor change in vacuum described by Eq. (15) or Eq. (19). In LMA-MSW, a $`\nu _e`$ gradually evolves into a $`\nu _2`$ within the sun, without anything sinusoidally oscillating. Then, this $`\nu _2`$ propagates to the earth without changing further, so that again nothing is oscillating. To be sure, a neutrino born as a $`\nu _e`$ has changed into a mixture of the flavors $`\nu _e`$ and $`\nu _x`$. But this flavor change has not involved any undulation.
### II.3 What Would We Like To Find Out?
Having reviewed what we have learned about the neutrinos from existing data, let us now ask what we would like to find out through future experiments. Let us consider just some of the most interesting open questions.
* How many neutrino species are there? Are there sterile neutrinos?
As has already been noted, we do not know how many neutrino mass eigenstates there are, with the LSND result suggesting that there are more than three. If there are indeed more than three mass eigenstates, then there are linear combinations of them that do not couple to the $`W`$ or $`Z`$ bosons. These linear combinations—so-called sterile neutrinos—do not participate in any of the known forces except gravity. Thus, confirmation of LSND would have far-reaching consquences. The MiniBooNE experiment, presently taking data, aims to confirm or refute LSND.
* What are the masses of the neutrino mass eigenstates?
This question has at least two parts. First, we know, within errors, the values of the atmospheric and solar splittings, $`\mathrm{\Delta }m_{\mathrm{atm}}^2`$ and $`\mathrm{\Delta }m_{\mathrm{sol}}^2`$. However, we do not know whether the solar pair, $`\nu _1`$ and $`\nu _2`$, separated by the smaller of the two splittings, $`\mathrm{\Delta }m_{\mathrm{sol}}^2`$, actually lies at the bottom of the spectrum, as in Fig. 3, or at the top. If this pair is at the bottom, the spectrum is called “normal”, and if it is at the top, the spectrum is called “inverted”. Whether the spectrum is normal or inverted is an interesting question, because generically, grand unified theories (GUTS) favor a normal spectrum r16 . This may be understood by remembering that GUTS relate the leptons to the quarks, and the quarks have normal spectra, with the most closely spaced quark pair at the bottom of the spectrum, not at the top. An inverted neutrino spectrum would be quite un-quark-like, and would probably involve a lepton symmetry, with no quark analogue, as the source of the near degeneracy of $`\nu _1`$ and $`\nu _2`$ at the top of the spectrum.
To determine whether the spectrum is normal or inverted, one may study the effect of earth matter on $`\nu _\mu \nu _e`$ and $`\overline{\nu _\mu }\overline{\nu _e}`$ oscillations. Using accelerator-generated $`\nu _\mu `$ and $`\overline{\nu _\mu }`$ beams, one will carry out a long-baseline experiment with an $`L/E`$ that makes it sensitive to the atmospheric mass gap, $`\mathrm{\Delta }m_{\mathrm{atm}}^2\mathrm{\Delta }m_{32}^2`$. The latter quantity is positive for a normal spectrum, but negative for an inverted one. By studying oscillations into a $`\nu _e`$ or $`\overline{\nu _e}`$, one will gain sensitivity to the $`W`$-exchange interaction betwen these particles and electrons in the earth matter through which the beam passes. As previously noted, this interaction leads to the extra interaction potential energy $`V_W=+\sqrt{2}G_FN_e`$ for a $`\nu _e`$ \[cf. Eq. (20)\], and to an equal but opposite extra energy for a $`\overline{\nu _e}`$. The signs of these extra energies are known from the Standard Model. By probing oscillations in which the quantity whose sign we wish to determine, $`\mathrm{\Delta }m_{32}^2`$, interferes with the $`W`$-exchange-induced extra $`\nu _e`$ or $`\overline{\nu _e}`$ energy, whose sign we already know, we can determine the sign of $`\mathrm{\Delta }m_{32}^2`$.
Since the experiment will involve $`\mathrm{\Delta }m_{\mathrm{atm}}^2`$ and focus on a $`\nu _e`$ or $`\overline{\nu _e}`$ final state, it will require that the $`\nu _e`$ fraction of $`\nu _3`$, the neutrino at one end of $`\mathrm{\Delta }m_{\mathrm{atm}}^2`$, be nonvanishing. Of course, this fraction is $`\mathrm{sin}^2\theta _{13}`$. Thus, not only CP violation in neutrino oscillation, but also one’s ability to tell whether the spectrum is normal or inverted by using matter effects, depends on a nonvanishing $`\theta _{13}`$ r17 .
In earth matter, the vacuum mixing angle $`\theta _{13}`$ is replaced by an effective mixing angle $`\theta _M`$ given by Eq. (36), with $`\theta `$ taken to be $`\theta _{13}`$. For a neutrino beam, the parameter $`x`$ in Eq. (36) is defined by Eq. (34), with $`\mathrm{\Delta }m_{}^2`$ now taken to be $`\mathrm{\Delta }m_{\mathrm{atm}}^2\mathrm{\Delta }m_{32}^2`$. For an antineutrino beam, $`x`$ reverses sign. Note that $`x`$ embodies the interference between $`\mathrm{\Delta }m_{32}^2`$ and $`V_W`$, enabling us to determine the unknown sign of the first from the known sign of the second, as just described.
From Eq. (39) and the experimental knowledge that $`\mathrm{cos}2\theta _{13}>0.91`$, we see that for neutrino beam energies $`E\stackrel{<}{}\mathrm{\hspace{0.33em}2}`$ GeV, we may expand the denominator of Eq. (36) to obtain
$$\mathrm{sin}^22\theta _M\mathrm{sin}^22\theta _{13}[1\text{+}\text{(}\text{}\text{)}𝒮\frac{E}{6\mathrm{GeV}}].$$
(56)
Here, $`𝒮`$ is the sign of $`\mathrm{\Delta }m_{32}^2`$ \[see Eq. (34)\], and the positive (negative) sign on the right-hand side of the equation is for a neutrino (antineutrino) beam.
At oscillation maximum, the $`\nu _\mu \nu _e`$ and $`\overline{\nu _\mu }\overline{\nu _e}`$ oscillation probabilities will be proportional to the $`\mathrm{sin}^22\theta _M`$ given by Eq. (56). Thus, at oscillation maximum,
$$\frac{P(\nu _\mu \nu _e)}{P(\overline{\nu _\mu }\overline{\nu _e})}\text{is}\{\begin{array}{c}>1;\text{Normal Spectrum}(𝒮=+1)\\ <1;\text{Inverted Spectrum}(𝒮=1)\end{array}.$$
As Eq. (56) shows, at, say, $`E`$ 2 GeV, this ratio can deviate from unity quite substantially.
A second aspect of the question “What are the masses of the neutrino mass eigenstates?” is the issue of the absolute scale of neutrino mass. How high above the zero of mass does the entire spectral pattern depicted in Figs. 2 and 3 lie? As already noted, neutrino flavor change, which depends only on squared-mass splittings, and not on individual masses, cannot answer this question. To be sure, flavor change does provide a lower bound on the mass of the heaviest mass eigenstate. Namely, since the mass of the lightest mass eigenstate is not below zero, the mass of the heaviest one cannot be less than $`\sqrt{\mathrm{\Delta }m_{\mathrm{atm}}^2}`$ 0.04 eV.
Approaches to going beyond a lower bound to an actual numerical value include the study of tritium $`\beta `$ decay, neutrinoless nuclear double $`\beta `$ decay, and cosmology. Cosmology has already provided an interesting upper bound. Cosmological data plus several cosmological assumptions imply that
$$\underset{i}{}m_i<(0.41.0)\mathrm{eV},$$
(57)
where $`m_i`$ is the mass of mass eigenstate $`\nu _i`$, and the sum runs over all the mass eigenstates r18 . Within the context of a three-neutrino spectrum with the observed squared-mass splittings, this bound implies that the mass of the heaviest $`\nu _i`$ is less than (0.2 – 0.4) eV. Thus, combining the flavor change and cosmological information,
$$0.04\mathrm{eV}<\text{Mass\hspace{0.17em}[Heaviest }\nu _i\text{]}<(0.20.4)\mathrm{eV}.$$
(58)
This is already an interesting indication of the scale of neutrino mass. However, we would like to learn more about the reliability of the cosmological assumptions, and we would like to aim for, not just bounds, but a numerical value.
* How large is $`\theta _{13}`$?
As we have seen, both CP violation and our ability to tell whether the mass spectrum is normal or inverted depend on $`\theta _{13}`$. Thus, we would like to know at least the approximate value of this angle. To see how we may learn that, let us recall that $`\mathrm{sin}^2\theta _{13}`$ is the small $`\nu _e`$ fraction of the isolated mass eigenstate $`\nu _3`$ \[see Fig. 3\]. This mass eigenstate is at one end of the atmospheric splitting, $`\mathrm{\Delta }m_{\mathrm{atm}}^2`$. Thus, to measure $`\theta _{13}`$, we must do an oscillation experiment whose $`L/E`$ makes it sensitive to $`\mathrm{\Delta }m_{\mathrm{atm}}^2`$, and that involves a $`\nu _e`$, either as the initial or final flavor. Two complementary approaches are being pursued. In the first, one will seek to observe the disappearance of reactor $`\overline{\nu _e}`$, which are particles of the required electron flavor, while they travel a distance $`L`$ 1.5 km. Since reactor $`\overline{\nu _e}`$ typically have an energy $`E`$ 3 MeV, $`L/E`$ will be $``$ 500 km/GeV, making the experiment sensitive, according to Eq. (16), to the atmospheric splitting $`\mathrm{\Delta }m_{\mathrm{atm}}^22.4\times 10^3`$ eV<sup>2</sup>. In the second approach, one will seek to observe the appearance of $`\nu _e`$, neutrinos of the required electron flavor, in an accelerator $`\nu _\mu `$ beam whose energy $`E`$ is of order 1 GeV, and which travels a distance $`L`$ of several hundred kilometers from source to detector. The $`L/E`$ of this experiment will be comparable to that of the reactor experiment.
It can easily be shown that the reactor experiment is sensitive to $`\theta _{13}`$, and $`\theta _{13}`$ alone, while the accelerator experiment is sensitive to $`\theta _{13}`$, $`\theta _{23}`$, the CP-violating phase $`\delta `$, and the normal or inverted character of the mass spectrum. Thus, the accelerator approach has the advantage of providing access to several neutrino properties that we wish to learn. However, these properties will have to be disentangled from one another in the data. This disentanglement will be facilitated by a clean measurement of $`\theta _{13}`$ from a reactor experiment.
* Are neutrinos their own antiparticles?
Alone among the fundamental fermionic constituents of matter, neutrinos might be their own antiparticles. A quark or charged lepton cannot be its own antiparticle, because it is electrically charged, and its antiparticle has the opposite electric charge. However, neutrinos are not electrically charged and, as far as we know, they do not carry any other conserved, charge-like quantity either. It might be thought that there is a conserved charge-like lepton number $`L`$, defined by
$$L(\nu )=L(\mathrm{})=L(\overline{\nu })=L(\overline{\mathrm{}})=1,$$
(59)
that distinguishes neutrinos $`\nu `$ and charged leptons $`\mathrm{}`$ from their antiparticles. However, there is absolutely no evidence that any such conserved quantum number exists. If it does not exist, then there is nothing to distinguish a neutrino from its antiparticle. Then each neutrino mass eigenstate $`\nu _i`$ is identical to its antiparticle $`\overline{\nu _i}`$, making the neutrinos very different from the charged leptons and the quarks.
A neutrino that is identical to its antiparticle is referred to as a Majorana particle, while one that is not is referred to as a Dirac particle.
Many theorists believe that, quite likely, the lepton number $`L`$ defined by Eq. (59) is not conserved, so that neutrinos are Majorana particles. One reason for this widely-held belief is the character of possible neutrino mass terms in extensions of the Standard Model (SM) r19 . The SM can be defined by the fields it contains, and by some general principles, such as weak isospin invariance and renormalizability. It is observed that anything allowed by the field content of the pre-neutrino-mass, original version of the SM, and allowed by the general SM principles, occurs in Nature.
The original version of the SM, which omits neutrino masses and contains no chirally right-handed neutrino fields, $`\nu _R`$, but only left-handed ones, $`\nu _L`$, conserves $`L`$. However, now that we know neutrinos have masses, we must extend the SM to accommodate them. Suppose that we try to do this in a way that will preserve the conservation of $`L`$. Then, for a neutrino $`\nu `$, we add to the SM Lagrangian a so-called “Dirac mass term” of the form
$$_D=m_D\overline{\nu _L}\nu _R+\mathrm{h}.\mathrm{c}..$$
(60)
Here, $`m_D`$ is a constant, and $`\nu _R`$ is a right-handed neutrino field that we are forced to add to the SM in order to construct the Dirac mass term. This mass term is of the same form as those that give masses to the charged leptons and quarks. It does not mix neutrinos and antineutrinos, and consequently conserves $`L`$.
In the SM, left-handed fermion fields belong to weak-isospin doublets, but right-handed ones are isospin singlets. Thus, the newly-introduced $`\nu _R`$ will carry no weak isospin. Hence, once $`\nu _R`$ is present, all the SM principles, notably including weak-isospin conservation, allow the occurrence of the “Majorana mass term”
$$_M=m_M\overline{\nu _R^c}\nu _R+\mathrm{h}.\mathrm{c}..$$
(61)
Here, $`m_M`$ is a constant, and $`\nu _R^c`$ is the charge conjugate of $`\nu _R`$. Since $`\nu _R`$ carries no weak isospin, neither does $`\nu _R^c`$, so $`_M`$ is fully weak-isospin conserving. However, any Majorana mass term of this form turns a $`\nu `$ into a $`\overline{\nu }`$, and a $`\overline{\nu }`$ into a $`\nu `$. Thus, it clearly does not conserve $`L`$.
The neutrino $`\nu `$ appearing in the mass terms of Eqs. (60) and (61) is not a mass eigenstate, but one of the underlying states in terms of which the model is written. Once the $`\nu \overline{\nu }`$ mixing, hence $`L`$-nonconserving, Majorana mass term of Eq. (61) is present, nothing can distinguish a neutrino mass eigenstate from its antiparticle, so that diagonalization of the model will yield mass eigenstates $`\nu _i`$ that are Majorana particles.
Nature has proved to contain (except perhaps in the Higgs sector) everything allowed by the general principles of the SM and the field content of its original version, which omits neutrino masses. It is then natural to expect that she also contains everything allowed by the principles of the SM and the field content of its extended version, which includes neutrino masses. If so, then Nature contains Majorana mass terms, so that $`L`$ is not conserved, and neutrinos are Majorana particles.
Confirmation that $`L`$ is not conserved would come from the observation of neutrinoless double beta decay ($`0\nu \beta \beta `$), the process Nucl $``$ Nucl$`{}_{}{}^{}+2e^{}`$, in which one nucleus decays to another with the emission of two electrons and nothing else. This process, with no leptons in the initial state but two in the final state, is manifestly lepton-number nonconserving, so its discovery would establish that neutrinos are Majorana particles.
It is expected that the process will be dominated by the diagram shown in Fig. 6. In this diagram, one or another of the neutrino mass eigenstates $`\nu _i`$ is exchanged between two virtual $`W`$ bosons to create the outgoing electrons. The $`0\nu \beta \beta `$ amplitude is then a coherent sum over the contributions of the different $`\nu _i`$.
It is assumed that the $`e\nu _iW`$ vertices in Fig. 6 are Standard Model vertices, involving the SM left-handed weak current. The exchanged $`\nu _i`$ produced in conjunction with an $`e^{}`$ by this current will have the helicity we normally associate with an antineutrino. That is, it will be dominantly right handed. However, just as the positron emitted in a $`\beta `$ decay is not 100% polarized, the virtual $`\nu _i`$ produced in the process of Fig. 6 will not be 100% polarized either. Rather, it will have a small component that is left handed. This component will be of order the mass $`m_i`$ of $`\nu _i`$, divided by its energy $`E`$. It is only this $`𝒪[m_i/E]`$ left-handed component that the left-handed current absorbing the exchanged $`\nu _i`$ can absorb without further suppression. Thus, the contribution of $`\nu _i`$ exchange to the $`0\nu \beta \beta `$ amplitude is proportional to $`m_i`$. It is also proportional to $`U_{ei}^2`$, since, as shown in Fig. 6, there is a factor of $`U_{ei}`$ at each of the $`e\nu _iW`$ vertices. Summing over the contributions of the different $`\nu _i`$, we conclude that the amplitude for $`0\nu \beta \beta `$ is proportional to the quantity
$$m_{\beta \beta }\left|\underset{i}{}m_iU_{ei}^2\right|.$$
(62)
This quantity is known as the effective Majorana neutrino mass for neutrinoless double beta decay.
That the $`0\nu \beta \beta `$ amplitude growing out of the diagram in Fig. 6 is proportional to neutrino mass is no surprise. After all, we have assumed that the $`e\nu _iW`$ vertices in Fig. 6 are SM vertices. But SM vertices conserve $`L`$. Thus, the $`L`$-nonconservation required by $`0\nu \beta \beta `$ must be coming from underlying Majorana neutrino mass terms. If the neutrino masses are turned off, the $`L`$\- nonconservation disappears.
It is easy to show that, regardless of what diagrams are actually involved in $`0\nu \beta \beta `$, the observation of this decay would still imply the existence of Majorana neutrino mass terms r20 . Thus, the observation of $`0\nu \beta \beta `$ would teach us that the physics underlying neutrino masses is of a different character than that underlying the masses of the quarks and charged leptons, none of which can have Majorana masses r21 .
* Do neutrino interactions violate CP? Is neutrino CP violation the reason we exist?
There are several compelling reasons to search for CP violation in neutrino interactions. First, we would like to know whether the leptonic interactions, like the quark interactions, violate CP invariance, or whether CP noninvariance is something peculiar to quarks. Secondly, we would like to know whether leptonic CP violation is responsible for the fact that the universe contains matter (of which we are made) but essentially no antimatter (which, if present, would annihilate us). That is, does leptonic CP violation explain why we exist?
Symmetry and other considerations suggest that, initially, the Big Bang produced equal amounts of matter and antimatter. Thus, the presently-observed preponderance of matter over antimatter must have developed after the initial instants. Now, if one starts with equal amounts of matter and antimatter, and the two behave identically, then one will continue to have equal amounts of the two. Thus, the development of an asymmetry—the present preponderance of matter over antimatter—requires that the two behave differently. That is, it requires a violation of CP invariance.
Experiments with $`K`$ and $`B`$ mesons have found CP violation in quark interactions. However, apart from some unconfirmed anomalies, this quark CP violation is very well described by the Standard Model, and it is known that SM quark CP violation is completely inadequate to explain the observed preponderance of matter over antimatter in the universe. Hence, it is very interesting to ask whether leptonic CP violation could explain it.
There is a very natural way in which leptonic CP violation could indeed explain the observed matter-antimatter asymmetry. The most popular theory of why neutrinos are so light is the See-Saw Mechanism r22 . This associates with each light neutrino $`\nu `$ a very heavy neutrino $`N`$. The heavier $`N`$ is, the lighter its see-saw partner, $`\nu `$, will be (hence the name “see-saw”). The masses of the heavy neutrinos $`N`$ are thought to be in the range $`10^{(915)}`$ GeV. Thus, these heavy neutrinos cannot be produced in laboratory experiments. However, like everything else, they would have been produced in the hot Big Bang. Now, it is a signature feature of the see-saw model that both the light neutrinos $`\nu `$ and their heavy see-saw partners $`N`$ are Majorana particles. Thus, an $`N`$ can decay both via $`N\mathrm{}+\mathrm{}`$ and via $`N\overline{\mathrm{}}+\mathrm{}`$, where $`\mathrm{}`$ is a charged lepton. However, if CP is violated in these CP-mirror-image leptonic decays, then
$$\mathrm{\Gamma }[N\mathrm{}+\mathrm{}]\mathrm{\Gamma }[N\overline{\mathrm{}}+\mathrm{}].$$
(63)
This rate inequality would have led to an early universe containing unequal numbers of leptons and antileptons. Non-perturbative Standard Model “sphaleron” processes would then have converted some of this lepton-antilepton asymmetry into a baryon-antibaryon asymmetry, resulting in the matter-antimatter asymmetric universe that we see today.
The production of the matter-antimatter asymmetry via CP violation in the decays of heavy neutrinos is known as Leptogenesis r23 . Obviously, we cannot confirm this scenario by repeating the early universe. However, we can lend credibility to it by demonstrating that CP is violated in the interactions of today’s light neutrinos $`\nu `$, which are the see-saw partners of the heavy neutrinos. We can demonstrate this light-neutrino CP violation by showing that the probabilities for the CP-mirror-image oscillations $`\nu _\alpha \nu _\beta `$ and $`\overline{\nu _\alpha }\overline{\nu _\beta }`$ are different.
It should be noted that even when, as in the see-saw model, the neutrino mass eigenstates are identical to their antiparticles, $`\nu _\alpha \nu _\beta `$ and “$`\overline{\nu _\alpha }\overline{\nu _\beta }`$” are still different processes. For example, when we make neutrinos via $`\pi ^+`$ decay, and look for the production of electrons by these neutrinos in a detector, we are studying the oscillation $`\nu _\mu \nu _e`$. However, when we make neutrinos via $`\pi ^{}`$ decay, and look for the production of positrons by these neutrinos, we are studying the oscillation conventionally referred to as “$`\overline{\nu _\mu }\overline{\nu _e}`$”. It is called $`\overline{\nu _\mu }\overline{\nu _e}`$ because it would involve antineutrinos if the latter differed from neutrinos. Regardless of notation, the oscillation involving a $`\pi ^+`$ at the source and an $`e^{}`$ at the detector, and its CP-mirror image, the oscillation involving a $`\pi ^{}`$ at the source and an $`e^+`$ at the detector, are clearly different processes. If the probabilities for these two processes, which are given by Eq. (15) even when the neutrino mass eigenstates are identical to their antiparticles r11 , differ, then CP is violated.
The relation between CP violation in (light) neutrino oscillation and in heavy $`N`$ decay is model dependent. However, if CP is violated in the oscillation of the light neutrinos, then quite likely it is also violated in the decays of the heavy neutrinos $`N`$ that are the see-saw partners of the light neutrinos r25 . Then leptogenesis stemming from $`N`$ decays may well have been the origin of the matter-antimatter asymmetry of the universe.
If $`N`$ decays did lead to the present proponderance of matter—of which we are made—over antimatter, then we are all descendants of heavy neutrinos.
## III CONCLUSION
Wonderful experiments have led to the discovery of neutrino mass. This discovery has raised some very interesting questions, and we would like very much to learn the answers to these questions.
###### Acknowledgements.
It is a pleasure to thank C. Albright, A. de Gouvêa, S. Parke, and L. Stodolsky for many useful conversations relevant to the topics of these lectures. Fermilab is operated by Universities Research Association Inc. under Contract No. DE-AC02-76CH03000 with the United States Department of Energy. |
warning/0506/physics0506024.html | ar5iv | text | # On the displacements of Einsteinian fields et cetera
## Abstract.
I give here: *i*) a very simple proof that the physical non-existence of gravitational waves (GW’s) is quite consistent with the basic principles of general relativity (GR); *ii*) a new argument against the physical existence of GW’s; *iii*) a criticism of Fock’s treatment of the GW’s; *iv*) some remarks on recent experimental investigations concerning the GW’s.
To be published on *Spacetime & Substance.*
email: angelo.loinger@mi.infn.it
Dipartimento di Fisica, Università di Milano, Via Celoria, 16 - 20133 Milano (Italy)
1. – The following is a widespread and erroneous opinion: *Without gravitational waves (GW’s), one would have to explain an* *instantaneous* *propagation of a change in the metric over the whole universe simply by changing the distribution of stress or mass of a given physical system*. – In reality, the physical non-existence of GW’s is quite consistent with the principles of general relativity (GR), as I have shown *ad abundantiam* in various papers , but perhaps in too concise ways insofar as the above specific belief is concerned. I shall give now in sects.2., 3. a detailed treatment of it, with the hope of convincing even the most naive among the physicists that the adjective “instantaneous” is not the attribute of a relativistic bugaboo – if it is properly understood.
In sect.4. I give a new argument against the physical existence of GW’s. In sect.5. Fock’s computations concerning the GW’s are critically examined. The *Appendix* reports some recent (negative) results of the experimental search of GW’s due to LIGO collaboration.
2. – In previous Notes I have repeatedly emphasized that Einstein field is *not* analogous to Maxwell field, since it has peculiar properties of its own that are not shared by the electromagnetic field. If, however, we neglect for a moment – *ad usum Delphini* – the existence of the e.m. waves, we can exploit a precise property of Maxwell field for our purpose. For convenience, I utilize here the treatment of Liénard-Wiechert e.m. fields – created by a moving point charge – as is developed in the well known treatise by Becker and Sauter ; see in particular p.293 of this book, which gives the expressions of the electric and magnetic fields due to Liénard and Wiechert. For our aim, it is expedient to consider the first part, say $`\text{E}_1`$, of the electric field E (e.g.), i.e. the part which does not depend on the charge acceleration. We have
(1)
$$\text{E}_1(\tau )/e=\left[\frac{(\text{r}r\text{v}/c)(1v^2/c^2)}{(r\text{r}\text{v}/c)^3}\right]_{\tau tr/c},$$
with evident and standard notations. As Becker and Sauter write, $`\text{E}_1`$ has the character of a *static* field, it falls off as $`1/r^2`$ for large distances. Since eq. (1) gives the first part of the expression of the global field E, which is valid for *all* velocities v, it must agree with the field, say $`\text{E}^{}`$, created by a uniformly moving charged particle (see sect.64 of ):
(2)
$$\text{E}^{}(t)/e=\frac{\text{r}(1v^2/c^2)}{[r^2(\text{r}\times \text{v}/c)^2]^{3/2}};$$
the formal difference between expressions (1) and (2) comes from the different meanings of the vector r in the two formulae. In eq.(2) $`\text{r}=\text{r}(t)`$ is set equal to the vector from the *instantaneous* particle location, say $`B`$, to the field point $`P`$, while in eq.(1) by $`\text{r}=\text{r}(\tau )\text{r}(tr/c)`$ we understand the vector from the particle location, say $`A`$, at *retarded* time $`\tau tr/c`$, to the field point $`P`$. For the case of *constant* velocity we obviously have:
(3)
$$\text{r}(t)=\text{r}(\tau )\frac{r(\tau )}{c}\text{v}.$$
If we write $`\text{r}(t)\text{r}_0`$, and $`\text{r}(\tau )\text{r}`$ (as in eq.(1)), we find from $`\text{r}_0=\text{r}r\text{v}/c`$ for the denominator of eq. (2) that
(4)
$$\left[r_0^2\left(\text{r}_0\times \text{v}/c\right)^2\right]^{3/2}=\left(r\text{r}\text{v}/c\right)^3,$$
i.e. the denominator of eq. (1). Thus the field $`\text{E}_1`$ actually represents the field *moving along with the particle*; and this is clearly true also for a *non*-constant speed. By contrast, the second part, $`\text{E}_2`$, of the total electric field $`\text{E}=\text{E}_1+\text{E}_2`$, which is proportional to the acceleration , has the character of a wavy $`(1/r)`$ \- decreasing field.
(I have reproduced almost literally some passages of Becker and Sauter , only the italics are mine.)
3. – We have seen that the *static* part $`\text{E}_1`$ of Liénard-Wiechert electric field E *moves* *en bloc* *with the particle*. Now, the *same* thing happens, in the *exact* formulation of GR, for the Einstein field $`g_{jk}(x^0,\text{x})`$ , $`(j,k=0,1,2,3)`$, since – as it has been proved *no “mechanism” exists in GR, which is capable of producing GW’s*. In other terms, if we displace a mass, its gravitational field and the related curvature of the interested manifold *displace themselves along with the mass*. In general, qualitatively speaking, we can affirm that under this respect Einstein field and Newton field behave in an identical way. This fact is mathematically and physically *evident* in Friedmann’s cosmological models, as I have shown , owing to the perfect agreement between Friedmann’s solutions and the solutions of corresponding Newtonian models. (Furthermore, we can remark that at any stage of the EIH-method of solution of field equations there is a suitable reference frame for which the solution has a *Newtonian* form.)
Conclusion: the widespread opinion reported at the beginning of sect.1 is *false*: the absence of GW’s does not generate any theoretical difficulty – as Levi-Civita had pointed out many years ago.
(Generally speaking, the real existence of *physical* waves requires the existence of *physically* privileged reference frames, or of a *material* medium as the cosmic ether. It is not the case of GR: in it a geodesic deviation must have a *Newton*-like character – and therefore could be recorded only by an apparatus in a relative proximity of the gravity source.) –
4. - The Einstein field equations share with Laplace-Poisson equation $`^2U=4\pi G\rho `$ an important property. Let us consider for a moment only the case of a “cloud of dust” with mass tensor $`T^{jk}=\rho u^ju^k`$, $`(j,k=0,1,2,3)`$, where $`\rho (x^0,\text{x})`$ is the invariant mass density and $`u^j(x^0,\text{x})`$ is the four-velocity of a gravitating particle. It is well known that we can always choose a Gaussian normal (“synchronous” in Landau’s terminology ) reference frame, for which:
(5)
$$\text{d}s^2=\left(\text{d}x^0\right)^2h_{\alpha \beta }(x^0,\text{x})\text{d}x^\alpha \text{d}x^\beta ,(\alpha ,\beta =1,2,3);$$
if there are only gravitational interactions – as in the present case –, this frame is also co-moving : the *world* lines of the “dust” particles are both *time* lines and *geodesic* lines. Our mass tensor $`T^{jk}`$ has only the component $`T^{00}=\rho `$ different from zero. *Thus* – exactly as it happens for Friedmann’s models *the metric tensor* $`g_{jk}(x^0,\text{x})`$ *depends* *only* *on* $`\rho (x^0,\text{x})`$, in perfect analogy with the Newtonian potential $`U`$, and it satisfies *identically* the geodesic equations. No GW’s are emitted – and this fact is now quite intuitive, because we see that the motion of the fluid has been formally “obliterated”.
This treatment can be immediately generalized to a continuum, whose particles are subject to gravitational and *non*-gravitational (e.g., electromagnetic) interactions. It is indeed sufficient to choose a *co-moving* reference frame – as it is always possible if the particle trajectories do not cross. Here too the metric tensor does not depend on the motion of the medium – motion that the metropolitan legend considers responsible of the emission of GW’s.
5. – Fock pretended erroneously that the so-called *harmonic frames* possess a *physical* privilege with respect to the other co-ordinate systems. Thus, in particular, all his computations concerning the GW’s are performed in a harmonic frame, and with mass tensors of *extended* bodies. Since the motions of gravitating *point* masses do not generate GW’s, it is difficult to believe in a thaumaturgical virtue of largeness. Indeed, the extended bodies are composed of particles, and, further, their translational motions are correctly treated as motions of material corpuscles.
Fock’s computations regarding the GW’s are rather poor in physical significance.
*APPENDIX*
*$`\alpha `$*) I report here the summary of a communication by I.Leonor at LIGO Scientific Collaboration meeting, March 23, 2005, entitled “Searching for GRB-GWB coincidence during LIGO science runs”.
*Summary*:
* developed scheme for searching for GRB-GWB coincidence in near real time
* looking forward to S5 run with $``$ 100 GRB triggers in one year of coincident run
* performed search for short-duration GW bursts coincident with S4, S3, and S2 GRB’s using crosscorrelation method
* sample probability distribution consistent with null hypothesis. –
The LIGO scholars are technically very clever, but evidently they cannot discover a non-existent object as a GW.
*$`\beta `$*) On *arXiv:gr-qc/0505029 v1* (6 May 2005) we can read a paper of 23 pages, written by 395 LIGO-researchers all over the world, entitled “Upper limits on gravitational wave bursts in LIGO’s second science run – LIGO-P040040-07-R”.
Here are some sentences from the ABSTRACT: “We perform a search for gravitational wave bursts using data from the second science run of the LIGO detectors, using a method based on a wavelet time-frequency decomposition. This search is sensitive to bursts of duration much less than a second and with frequency content in the 100-1100 Hz range. It features significant improvements in the instrument sensitivity and in the analysis pipeline with respect to the burst search previously reported by LIGO. $`[\mathrm{}]`$. No gravitational wave signals were detected in 9.98 days of analyzed data. $`[\mathrm{}]`$”.
At p.11 we read: “The WaveBurst analysis applied to the S2 data yielded 16 coincidence events (at zero-lag). The application of the *r*-statistic cut rejected 15 of them, leaving us with a single event that passed all the analysis criteria.”. And at p.13: “The investigation revealed that the event occurred during a period of strongly elevated acoustic noise at Hanford lasting tens of seconds, as measured by microphones placed near the interferometers. $`[\mathrm{}]`$. The source of the acoustic noise appears to have been an aircraft.” –
In spite of the repeated failures, the LIGO scientists are still hopeful. *Spes ultima dea*. –
*$`\gamma `$*) On *arXiv:gr-qc/0505042 v1* (10 May 2005) the above 395 scholars have published an article (7pp.) entitled “Search for Gravitational Waves from Primordial Black Hole Binary Coalescences in the Galactic Halo”. From the ABSTRACT: “We use data from the second science run of the LIGO gravitational-wave detectors to search for the gravitational waves from primordial black hole (PBH) binary coalescence with component masses in the range 0.2–$`1.0M_{}`$. $`[\mathrm{}]`$. No inspiral signals were found.” – Obviously: both GW’s and BH’s are non-existent objects . The so-called *observed* BH’s are enormously massive bodies restricted in relatively small volumes – as it can be demonstrated by a careful scrutiny of the concerned papers . –
*$`\delta `$*) Again the mentioned 395 scientists on *arXiv:gr-qc/0505041 v1* (12 May 2005): “Search for gravitational waves from galactic and extra–galactic binary neutron stars” (20pp.). From the ABSTRACT: “We use 373 hours ($``$ 15 days) of data from the second science run of the LIGO gravitational-wave detectors to search for signals from binary neutron star coalescences within a maximum distance of about 1.5 Mpc, a volume of space which includes the Andromeda Galaxy and other galaxies of the Local Group of galaxies. $`[\mathrm{}]`$. No inspiral gravitational wave events were identified in our search.” The conclusion of the paper is the following (p.19): “In this paper, we have presented a data analysis strategy that could lead to a detection of gravitational waves from binary neutron star inspirals. The methods used to validate the search illustrate the subtleties of the analysis of several detectors with different sensitivities and orientations. Moreover, the experience gained by following up the largest coincident triggers will be crucial input to investigations of event candidates that are identified in future searches.” An Italian jest says: *Chi vive sperando muore cantando*. – |
warning/0506/math0506094.html | ar5iv | text | # A note on Bruhat decomposition of 𝐺𝐿(𝑛) over local principal ideal rings
## 1. Introduction
Let $`A`$ be a local principal ideal commutative ring and let $`\mathrm{}=(\pi )`$ denote its maximal ideal. Denote by $`k`$ the length of the ring, that is, the least $`k`$ such that $`\mathrm{}^k=0`$ ($`k`$ might be infinite). Let $`B`$ denote the subgroup of upper triangular matrices in $`G=GL_n(A)`$, the group of invertible matrices with entries in $`A`$. This paper concerns the description of the double coset space $`B\backslash G/B`$. Since $`B`$ is the stabiliser of the standard flag in $`A^n`$, this space corresponds to the possible relative positions of two flags that are isomorphic to the standard flag (these are the flags whose reductions modulo $`\mathrm{}`$ are full flags of vector spaces in $`(A/\mathrm{})^n`$). We refer to such a flag as a *full free primitive flag over $`A`$ in $`A^n`$*. If
$``$ $`=\text{the space of full free primitive flags over }A\text{ in }A^n`$
$`B`$ $`=\mathrm{Stab}_G(\text{standard flag})`$
then one has
$$\times _GB\backslash G/B.$$
A complete description of the latter space for all $`n`$ contains the embedding problem of pairs of $`A`$-modules, which is known to be of *wild type* in general \[RS, Sch04\]. Thus, one does not expect a reasonably closed solution, and we aim at the more modest goal of constructing some invariants of the double coset space, describing the relations between them and exploring to what extent they can distinguish double cosets. When $`k=1`$, in which case the ring $`A`$ is a field, Bruhat decomposition says that $`B\backslash G/B`$ is parameterised by the symmetric group. In particular, the number of double cosets does not depend on the field. The first natural question is: for which pairs $`(n,k)`$ is the parametrisation independent of the ring? The answer to this question is given in Theorem 6.1. The first two instances in which a dependence on the ring occurs are $`(3,k)`$ with $`k>2`$ and $`(4,2)`$. For $`(3,k)`$ we give a complete parametrisation of the double coset space and give estimates for its size when $`A`$ is a finite ring (Section 4).
### 1.1. Related problems
Let $`P_1,P_2<G`$ be finite groups. Let $`\rho _i=\mathrm{Ind}_{P_i}^G1=[G/P_i]`$ be the representation of $`G`$ induced from the trivial representation of $`P_i`$ over $``$. The module of intertwining operators $`\mathrm{Hom}_G(\rho _1,\rho _2)`$ can be identified with the subalgebra $`[P_2\backslash G/P_1]`$ of left-$`(P_2,P_1)`$-invariant elements in the group algebra $`[G]`$ via the map $`[P_2\backslash G/P_1]\mathrm{Hom}(\rho _1,\rho _2)`$ given by
$$fT_f\text{ for each }f[P_2\backslash G/P_1],$$
where
$$T_fh(x)=\underset{gG}{}h(xg)f(g)\text{ for each }h\rho _1=[G/P_1].$$
Let $`A_i=A/\mathrm{}^i`$ for $`i`$ be the inverse system of the finite length quotients. Isomorphism types of finitely generated $`A_i`$-modules correspond to Young diagrams with height bounded by $`i`$. The Young diagram given by $`\lambda =(\lambda _1,\mathrm{},\lambda _j)`$, with $`i\lambda _1\mathrm{}\lambda _j0`$ corresponds to the $`A`$-module
$$M_\lambda =_{r=1}^jA/\mathrm{}^{\lambda _r}.$$
Let $`G_i=GL_n(A_i)`$, $`P_\lambda `$ be the stabiliser of a submodule of type $`\lambda `$ in $`A_i^n`$, and $`B_i`$ be the stabiliser of a full flag of free submodules in $`A_i^n`$.
Assume now that the residue field of $`A`$ is finite of order $`q`$. The induced representations $`\rho _i=\mathrm{Ind}_{B_i}^{G_i}1`$ play a significant role in the representation theory of the groups $`G_i`$ \[Hil93, Hil95\] in analogy with the role played by $`\rho _1`$ in the representation theory of $`GL_n`$ over finite fields \[Zel81\]. The latter representation is studied in terms of the Hecke algebra $`_{A,1}`$, where
$$_{A,i}=\text{End}_{G_i}(\rho _i)[B_i\backslash G_i/B_i]$$
and one has $`_{A,1}[S_n]`$, independent of the characteristic of the residue field. The algebra $`_{A,i}`$ continues to play an important role for $`i>1`$, however, its structure depends on the characteristic of the residue field. As a starting point, one would like to know its dimension - hence parameterise the double coset space $`B_i\backslash G_i/B_i`$. The number of double cosets depends on the residue field unless $`n2`$ or $`i=1`$ or $`n=3`$ and $`i=2`$ (see Theorem 6.1).
Broadening the frame a bit, it is natural to consider the category of diagrams over finite length $`A`$-modules; For a quiver $`Q`$ let $`𝒟_Q=\text{Fun}(Q,A\text{-mod})`$ be the category of functors from $`Q`$ (considered as a category) to the category of finitely generated $`A`$-modules. These categories occur naturally in the study of the above representations, in particular as the underlying sets for the modules of intertwining operators between representations. For example \[BO\]
$`P_{i^m}\backslash G_i/P_{i^m}`$ $`\{\text{isomorphism types of submodules of }M_{i^m}=A_i^m\}\text{Isom}(𝒟_{\{\}})`$
$`P_{i^m}\backslash G_i/P_\lambda `$ $`\{\text{isomorphism types of pairs }(NN^{})M_\lambda \}\text{Isom}(𝒟_{\{\}})`$
$`P_{i^m}\backslash G_i/B_i`$ $`\{\text{isomorphism types of chains submodules of }A_i^m\}\text{Isom}(𝒟_{\{\mathrm{}\}})`$
The category $`𝒟_{\{\}}`$ is nothing but the category of $`A`$-modules. The full subcategory of $`𝒟_{\{\}}`$ which consists of embeddings is of wild type \[RS, Sch04\]. The full subcategory of $`𝒟_{\{\mathrm{}\}}`$ which consists of embeddings is discussed in \[Sim02\].
### 1.2. Notations
Throughout this paper $`\pi `$ denotes a generator of the maximal ideal $`\mathrm{}`$ in $`A`$ and the order of $`\pi `$ is denoted by $`k`$, and might be infinite (as in the preceding paragraphs). The valuation of a non-zero element $`xA`$ is denoted by $`v(x)`$. For convenience we also write $`v(0)=k`$. For any ideal $`I`$ in $`A`$, $`v(I)`$ is the semigroup of valuation values of elements in $`I`$. The residue field of $`A/\mathrm{}`$ is denoted by $`𝐤`$, $`A^\times `$ denotes the multiplicative group, and $`A_i=A/\mathrm{}^i`$.
### 1.3. Acknowledgements
We thank the Tata Institute of Fundamental Research, in particular Ravi Rao and Dipendra Prasad, for bringing us together and for their warm hospitality, which resulted in this manuscript. We thank Markus Schmidmeier for supplying detailed information on the embedding problem. The first author also thanks Uri Bader for discussions which partly motivated this work.
## 2. Invariants
### 2.1. Upper triangular row and column operations
The determination of the double coset space $`B\backslash G/B`$ is part of the more general question of determining the orbits of the action of invertible upper triangular matrices by left and right multiplications on the set $`M_{nm}`$ of $`n\times m`$ matrices with entries in $`A`$. When two matrices $`\alpha `$ and $`\alpha ^{}`$ lie in the same orbit, we write $`\alpha \alpha ^{}`$. Let $`R_i`$ denote the $`i`$th row and $`C_j`$ denote the $`j`$th column of a matrix. Then two matrices lie in the same orbit if one can be obtained from the other by the following types of row and column operations:
* Multiplication of rows or columns by scalars (action of the ‘torus’)
$$R_iaR_i,\text{ with }aA^\times \text{ and }C_jaC_j,\text{ with }aA^\times .$$
* Addition of certain rows/columns to others (action of the ‘unipotent subgroup’)
$`R_iR_i+{\displaystyle \underset{i^{}>i}{}}a_i^{}R_i^{}\text{ with }a_i^{}A`$
$`C_jC_j+{\displaystyle \underset{j^{}<j}{}}a_j^{}C_j^{}\text{ with }a_j^{}A`$
Note that only row (respectively, column) operations with $`i^{}>i`$ (respectively, $`j^{}<j`$) are allowed.
Since row operations commute with column operations and scaling operations normalise addition operations, two matrices in $`M_{nm}`$ lie in the same orbit if and only if one can be obtained from the other by a sequence of scaling operations followed by row operations $`R_iR_i+_{i^{}>i}a_i^{}R_i^{}`$ with $`i`$ increasing from $`1`$ to $`n`$ followed by column operations $`C_jC_j+_{j^{}<j}a_j^{}C_j^{}`$ with $`j`$ decreasing from $`n`$ to $`1`$.
### 2.2. Decomposability
We discuss a class of matrices for which the problem of determining whether two matrices lie in the same double coset reduces to similar problems involving smaller matrices.
###### Proposition 2.1.
Suppose that $`n=n_1+n_2`$ and the the matrices $`\alpha `$ and $`\alpha ^{}`$ have block matrix decompositions
$$\alpha =\left(\begin{array}{cc}0& \alpha _1\\ \alpha _2& 0\end{array}\right)\text{ and }\alpha ^{}=\left(\begin{array}{cc}0& \alpha _1^{}\\ \alpha _2^{}& 0\end{array}\right),$$
with $`\alpha _iGL_{n_i}(A)`$. Then $`\alpha \alpha ^{}`$ if and only if $`\alpha _1\alpha _1^{}`$ and $`\alpha _2\alpha _2^{}`$.
###### Proof.
Write
$$\left(\begin{array}{cc}b_1& X\\ 0& b_2\end{array}\right)\left(\begin{array}{cc}0& \alpha _1\\ \alpha _2& 0\end{array}\right)\left(\begin{array}{cc}c_1& Y\\ 0& c_2\end{array}\right)=\left(\begin{array}{cc}0& \alpha _1^{}\\ \alpha _2^{}& 0\end{array}\right)$$
where $`b_1`$ and $`c_2`$ are upper-triangular invertible $`n_1\times n_1`$-matrices, $`b_2`$ and $`c_1`$ are upper triangular invertible $`n_2\times n_2`$ matrices, $`X`$ is an $`n_1\times n_2`$ matrix and $`Y`$ is an $`n_2\times n_1`$ matrix. Multiplying out, and comparing the the lower-right blocks gives $`b_2\alpha _2Y=0`$, which, since $`b_2`$ and $`\alpha _2`$ are invertible, implies that $`Y=0`$. Equating the remaining entries, and setting $`Y=0`$ gives that $`b_2\alpha _2c_1=\alpha _2^{}`$, $`b_1\alpha _1c_2=\alpha _1^{}`$ and $`X=b_1\alpha _1\alpha _2^1`$. This shows that $`\alpha \alpha ^{}`$ if and only if $`\alpha _1\alpha _1^{}`$ and $`\alpha _2\alpha _2^{}`$. ∎
This proposition shows that the classification of the double cosets for a given $`n`$ implies the classification for all $`n^{}<n`$. The following corollary allows one to reduce the equivalence problem to smaller $`n`$ for many matrices:
###### Corollary 2.2.
Suppose that $`n=n_1+n_2`$ and the the matrices $`\alpha `$ and $`\alpha ^{}`$ have block matrix decompositions
$$\alpha =\left(\begin{array}{cc}X_1& \alpha _1\\ \alpha _2& X_2\end{array}\right)\text{ and }\alpha ^{}=\left(\begin{array}{cc}X_1^{}& \alpha _1^{}\\ \alpha _2^{}& X_2^{}\end{array}\right),$$
with $`\alpha _iGL_{n_i}(A)`$. Then $`\alpha \alpha ^{}`$ if and only if $`\alpha _1X_1\alpha _2^1X_2\alpha _1^{}X_2^{}\alpha _2^1X_2^{}`$ and $`\alpha _2\alpha _2^{}`$.
###### Proof.
The upper triangular row and column operations described in Section 2.1 can be used to reduce $`\alpha `$ and $`\alpha ^{}`$ to matrices of the type that occur in Proposition 2.1, but with $`\alpha _1`$ and $`\alpha _1^{}`$ being replaced by $`\alpha _1X_1\alpha _2^1X_2`$ and $`\alpha _1^{}X_2^{}\alpha _2^1X_2^{}`$ respectively. ∎
### 2.3. Intersection Invariants
Let $`𝙵_0`$ denote the *standard flag* in $`A^n`$:
$$𝙵_0=(0=𝙵_0^0𝙵_0^1\mathrm{}𝙵_0^n=A^n)$$
where $`𝙵_0^i`$ is the $`A`$-module spanned by $`\{𝐞_1,\mathrm{},𝐞_i\}`$, and $`𝐞_i`$ is the $`i`$th standard basis vector in $`A^n`$. $`G`$ acts transitively on the set of full free primitive flags over $`A`$ in $`A^n`$. Thus the space of such flags is identified with $`G/B`$. For $`\alpha GL_n(A)`$, consider the corresponding flag $`𝙵=\alpha 𝙵_0`$ given by
$$𝙵=(0=𝙵^0\mathrm{}𝙵^n=A^n)=\alpha 𝙵_0\text{ where }𝙵^i=\alpha 𝙵_0^i.$$
Clearly the isomorphism classes of the $`A`$-modules $`𝙵^j𝙵_0^i`$ are invariants of the double coset. We will call these the intersection types. The intersection types are related to the column spaces of lower-left submatrices of $`\alpha `$ as follows: let $`[\alpha ]^{ij}`$ denote the lower-left $`(ni)\times j`$ submatrix.
###### Proposition 2.3.
The column space of $`[\alpha ]^{ij}`$ is isomorphic to $`𝙵^j/𝙵^j𝙵_0^i`$ as an $`A`$-module.
###### Proof.
The map from $`𝙵^j`$ to the column space of $`[\alpha ]^{ij}`$ is defined by taking the last $`ni`$ entries of a vector. The kernel is clearly $`𝙵^j𝙵_0^i`$. ∎
Furthermore, $`𝙵`$ induces a filtration on each graded piece $`𝙵_0^i/𝙵_0^{i1}`$ of the standard flag $`𝙵_0`$:
$$0=\frac{𝙵^0𝙵_0^i}{𝙵^0𝙵_0^{i1}}\frac{𝙵^1𝙵_0^i}{𝙵^1𝙵_0^{i1}}\mathrm{}\frac{𝙵^n𝙵_0^i}{𝙵^n𝙵_0^{i1}}=𝙵_0^i/𝙵_0^{i1}.$$
The $`j`$th graded piece of the above filtration is:
(1)
$$\frac{𝙵^j𝙵_0^i}{𝙵^{j1}𝙵_0^i+𝙵^j𝙵_0^{i1}}.$$
Being a subquotient of $`𝙵_0^i/𝙵_0^{i1}A`$, it must be isomorphic to $`A/(p^{r_{ij}})`$ for some $`0r_{ij}k`$ (if $`k=\mathrm{}`$, then some of the $`r_{ij}`$’s will be infinite). The $`B`$-action on the space $`G/B`$ of flags preserves the isomorphism classes of the $`A`$-modules in (1). Consider the matrix $`r(\alpha )=(r_{ij})`$. The above considerations show that it is invariant under left and right multiplications in $`B`$, and that each column sums to $`k`$. A similar argument can be used to show that each row sums to $`k`$. We call $`r(\alpha )`$ *the matrix of intersection numbers* of $`\alpha `$. When $`k=1`$, the matrix of intersection numbers is a permutation matrix, and is in fact, the unique permutation matrix that lies in the double coset of $`\alpha `$. *In this sense, the matrix of intersection numbers is a direct generalisation of the permutation associated to a matrix over a field by the Bruhat decomposition.*
### 2.4. Permutation
Consider the surjection induced by reduction modulo $`\mathrm{}`$
$$𝐖:B(A)\backslash GL_n(A)/B(A)B(𝐤)\backslash GL_n(𝐤)/B(𝐤)S_n$$
and view the double cosets as fibres over the field case. The image of $`\alpha G`$ in $`B(𝐤)\backslash GL_n(𝐤)/B(𝐤)`$ determines an $`n\times n`$ permutation matrix $`𝐖(\alpha )`$ by the Bruhat decomposition which is an invariant of the double coset of $`\alpha `$. Given a permutation matrix $`w`$ of order $`n`$, let $`N(w)`$ denote the number of double cosets for which $`𝐖=w`$. Say that the permutation $`w`$ is decomposable if there exists a partition $`n=n_1+n_2`$ and $`w_1`$ and $`w_2`$ permutation matrices of order $`n_1`$ and $`n_2`$ respectively such that
$$w=\left(\begin{array}{cc}0& w_1\\ w_2& 0\end{array}\right)$$
as a block matrix. By Corollary 2.2, we have:
(2)
$$N(w)N(w_2)\text{ for all }k.$$
## 3. The $`n=2`$ case
In this case there are $`k+1`$ double cosets. Geometrically, these double cosets parameterise the possible intersections of two free primitive sub-modules of $`A^2`$ of rank $`1`$. The intersection of two such submodules is a submodule of each of them, and hence isomorphic to $`A/\mathrm{}^r`$ for some $`0rk`$. We see that the intersection diagram, the intersection types and the intersection numbers carry the same information; in fact they are complete invariants. In terms of the permutation invariants, the fibre over the trivial permutation consists of $`k`$ elements corresponding to $`r=1,\mathrm{},k`$, and the fibre over the non-trivial permutation consists of one element corresponding to $`r=0`$. In terms of matrices, the set
$$\left\{\left(\begin{array}{cc}1& 0\\ \pi ^r& 1\end{array}\right)\right|0rk\}$$
is a complete set of representatives. The intersection numbers are given by
$$\alpha B\left(\begin{array}{cc}1& 0\\ \pi ^r& 1\end{array}\right)B\text{ if and only if }𝐫(\alpha )=\left(\begin{array}{cc}r& kr\\ kr& r\end{array}\right).$$
## 4. The $`n=3`$ case
### 4.1. Fibration over the residue field
We have the following description of the fibres over permutations:
(3)
$$\begin{array}{ccccccc}B\backslash GL_3(A)/B& B\backslash M_2^{}/B& v(\mathrm{})\times v(\mathrm{})& v(\mathrm{})\times v(\mathrm{})& v(\mathrm{})& v(\mathrm{})& \{\}\\ & & & & & & \\ B\backslash GL_3(𝐤)/B& \left(\begin{array}{ccc}1& 0& 0\\ 0& 1& 0\\ 0& 0& 1\end{array}\right)& \left(\begin{array}{ccc}0& 1& 0\\ 1& 0& 0\\ 0& 0& 1\end{array}\right)& \left(\begin{array}{ccc}1& 0& 0\\ 0& 0& 1\\ 0& 1& 0\end{array}\right)& \left(\begin{array}{ccc}0& 0& 1\\ 1& 0& 0\\ 0& 1& 0\end{array}\right)& \left(\begin{array}{ccc}0& 1& 0\\ 0& 0& 1\\ 1& 0& 0\end{array}\right)& \left(\begin{array}{ccc}0& 0& 1\\ 0& 1& 0\\ 1& 0& 0\end{array}\right)\\ S_3& 1& s_1& s_2& s_1s_2& s_2s_1& s_1s_2s_1\end{array}$$
Except for the fibre over the trivial element (for which the notation used in the table is explained below), one easily verifies that the fibres are indeed the ones written above. However, perhaps for the case of the double cosets lying over the permutations labelled $`s_1`$ and $`s_2`$ a remark is in order: any element lying above $`s_1`$ can be brought to the form
$$\left(\begin{array}{ccc}0& 1& 0\\ 1& 0& 0\\ \pi ^i& \pi ^j& 1\end{array}\right),1i,jk$$
We see that these lie in different double cosets by observing that they have different intersection types. A similar argument holds for $`s_2`$.
As for the fibre over $`\{1\}`$, it is determined by the $`2\times 2`$ lower left sub-matrix. We are therefore led to analyse the double coset space of $`2\times 2`$ matrices $`B\backslash M_2^{}/B`$. Here $`M_2^{}`$ denotes the set of those $`2\times 2`$ matrices for which only the top right entry is a unit and (with a slight abuse of notation) $`B`$ denotes the group of upper triangular matrices in $`GL_2(A)`$.
### 4.2. The space $`B\backslash M_2^{}/B`$
For $`\alpha M_2^{}`$, let $`v(\alpha )`$ denote the valuation matrix. The matrix is in standard form if the valuations of the non-zero entries form a standard tableaux (decreasing with respect to column numbers and increasing with respect to row numbers). Every matrix in $`M_2^{}`$ can be reduced to standard form. Assume that $`\alpha `$ is in standard form. Write
$$v(\alpha )=\left(\begin{array}{cc}i& 0\\ j& l\end{array}\right).$$
If at least one of the entries of $`i`$, $`j`$ or $`l`$ is $`k`$, the double coset is completely determined by the valuation matrix, and the classification of the orbits is given completely in terms of the matrix of intersection numbers.
If none of the entries of $`v(\alpha )`$ is equal to $`k`$, the matrix could be brought to the form:
$$\alpha (a)=\left(\begin{array}{cc}\pi ^i& 1\\ \pi ^ja& \pi ^l\end{array}\right)$$
with $`k>j>\mathrm{max}\{i,l\}\mathrm{min}\{i,l\}>0`$ and $`aA^\times `$ (in fact can be taken in $`A_{kj}^\times `$).
We shall consider the cases where $`k`$ is finite and infinite simultaneously. When $`k`$ is finite, our equations are modulo a power of $`\pi `$. When $`k`$ is infinite, then $`A`$ is a domain, and the process of going modulo powers of $`\pi `$ should be ignored. Let $`ϵ=\mathrm{min}\{ji,jl,i,l\}`$ and $`\delta (a)=v(a1)`$.
###### Proposition 4.1.
Let $`i,j`$ and $`l`$ satisfy $`k>j>\mathrm{max}\{i,l\}\mathrm{min}\{i,l\}>0`$. For every $`a,a^{}A_{kj}^\times `$, let $`\alpha (a)`$ and $`0\delta (a)kj`$ be as above. We have
1. If $`ji+l`$, then
$$\alpha (a)\alpha (a^{})aa^{}mod\pi ^{\mathrm{min}\{ϵ,kj\}}$$
2. If $`j=i+l`$, then
$$\alpha (a)\alpha (a^{})\{\begin{array}{cc}\delta (a)=\delta (a^{})=:\delta \hfill & \\ aa^{}mod\pi ^{\mathrm{min}\{ϵ+\delta ,kj\}}\hfill & \end{array}$$
###### Proof.
$`\alpha (a)\alpha (a^{})`$ if and only if there exist $`x_{11},x_{22},y_{11},y_{22}A^\times `$ and $`x_{12},y_{12}A`$ such that
(4)
$$\left(\begin{array}{cc}x_{11}& x_{12}\\ 0& x_{22}\end{array}\right)\left(\begin{array}{cc}\pi ^i& 1\\ \pi ^ja& \pi ^l\end{array}\right)=\left(\begin{array}{cc}\pi ^i& 1\\ \pi ^ja^{}& \pi ^l\end{array}\right)\left(\begin{array}{cc}y_{11}& y_{12}\\ 0& y_{22}\end{array}\right)$$
Equating the entries on both sides of (4) we get the following system of equations:
$`\pi ^ix_{11}+\pi ^jax_{12}`$ $`=`$ $`\pi ^iy_{11}`$
$`x_{11}+\pi ^lx_{12}`$ $`=`$ $`\pi ^iy_{12}+y_{22}`$
$`\pi ^jax_{22}`$ $`=`$ $`\pi ^ja^{}y_{11}`$
$`\pi ^lx_{22}`$ $`=`$ $`\pi ^ja^{}y_{12}+\pi ^ly_{22}`$
and after successive substitutions of these equations we obtain
(5)
$$(1a^{}a^1)y_{11}=(\pi ^{ji}a+\pi ^l)x_{12}+(\pi ^i\pi ^{jl}a^{})y_{12}mod\pi ^{kj}.$$
The equation (5) can be solved if and only if the coefficient of $`y_{11}`$ on the left hand side lies in the ideal generated by the coefficients of $`x_{12}`$ and $`y_{12}`$ on the right hand side, in other words, if and only if
(6)
$$v(aa^{})\mathrm{min}\{v(\pi ^{ji}a+\pi ^l),v(\pi ^ia^{}\pi ^{jl}),kj\}.$$
If $`ji+l`$, then the right hand side of (6) is $`\mathrm{min}\{ϵ,kj\}`$. This proves the first part of the proposition. To prove the second part, we observe that if $`j=i+l`$, then (6) becomes
(7)
$$v(aa^{})\mathrm{min}\{l+v(a1),i+v(a^{}1),kj\}.$$
However, equating the valuations of the determinants of both sides of (4) gives
$$\mathrm{min}\{kj,v(a1)\}=\mathrm{min}\{kj,v(a^{}1)\},$$
and taking this condition into account allows us to replace $`a^{}`$ with $`a`$ in the right hand side of (6) to get $`\mathrm{min}\{kj,i+v(a1),l+v(a1)\}`$, which, when $`j=i+l`$, is the same as $`\mathrm{min}\{kj,ϵ+\delta \}`$.
Conversely, it is easy to see that a solution of (5) can always be extended to a solution of (4). ∎
### 4.3. Parametrisation of the double coset space
###### Corollary 4.2.
The double coset space $`B\backslash GL_3(A)/B`$ is parameterised by
$$B\backslash GL_3(A)/B=(B\backslash M_2^{}/B)(v(\mathrm{})\times v(\mathrm{}))(v(\mathrm{})\times v(\mathrm{}))v(\mathrm{})v(\mathrm{})\{\}$$
where
$$\begin{array}{c}\hfill B\backslash M_2^{}/B=\underset{j=2}{\overset{k1}{}}\left[\left(\underset{\begin{array}{c}1i,lj1\\ i+lj\end{array}}{}A_{\mathrm{min}\{ϵ,kj\}}^\times \right)\left(\underset{\begin{array}{c}1i,lj1\\ i+l=j\end{array}}{}\underset{\delta =0}{\overset{kj}{}}A_{\mathrm{min}\{ϵ+\delta ,kj\}}^{\times ,\delta }\right)\right]\\ \hfill \left\{\left(\begin{array}{cc}i& 0\\ j& l\end{array}\right)\right|[j=k,1i,lk]\text{or}[i<j<k=l]\text{or}[l<j<k=i]\}\end{array}$$
and $`A_i^{\times ,\delta }=\{aA_i^\times |v(a1)=\delta \}`$.
###### Corollary 4.3.
If the residue field of $`A`$ is finite with $`q`$ elements then
$$p_1(k)q^{k/3}<|B\backslash GL_3(A_k)/B|<p_2(k)q^{k/3}$$
for some positive polynomials $`p_1`$ and $`p_2`$.
###### Proof.
For each permutation except for the trivial one, the number of double cosets that lie over that permutation has polynomial growth in $`k`$, hence does not affect the bounds. This is still the case for the identity permutation, when at least one of $`i`$, $`l`$ or $`j`$ is $`k`$, as there is only one double coset of the standard form of §4.2 with valuations $`i`$, $`l`$ and $`j`$.
Otherwise, different powers of $`q`$ appear as the number of double cosets with fixed values of $`i`$, $`l`$ and $`j`$. Here, there are two possible cases:
$`ji+l`$ \# cosets $`=|A_{\mathrm{min}\{ϵ,kj\}}^\times |=q^{\mathrm{min}\{ϵ,kj\}}(1q^1)`$
$`j=i+l`$ \# cosets $`=|{\displaystyle \underset{\delta =0}{\overset{kj}{}}}A_{ϵ+\delta ,kj}^{\times ,\delta }|={\displaystyle \underset{\delta =0}{\overset{kj}{}}}q^{\mathrm{min}\{ϵ,kj\delta \}}`$
Clearly, the highest power of $`q`$ that can occur in this way is $`k/3`$, hence the upper bound. This value is indeed achieved when $`i`$, $`l`$ and $`j/2`$ are close to $`k/3`$, hence the lower bound. ∎
## 5. The $`n=4`$, $`k=2`$ case
We saw that when $`n=3`$ the number of double cosets does not depend on the ring (e.g., the characteristic of the residue field) when $`k=2`$. This is no longer true when $`n=4`$. We illustrate this by describing some of the double cosets lying over the permutation
$$w=\left(\begin{array}{cccc}0& 1& 0& 0\\ 1& 0& 0& 0\\ 0& 0& 0& 1\\ 0& 0& 1& 0\end{array}\right)$$
Let $`\sigma =\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)`$, and $`\tau (a)=\left(\begin{array}{cc}\pi & a\pi \\ \pi & \pi \end{array}\right)`$.
###### Proposition 5.1.
For $`a,a^{}A`$, the (block) matrices
$$\left(\begin{array}{cc}\sigma & 0\\ \tau (a)& \sigma \end{array}\right)\text{ and }\left(\begin{array}{cc}\sigma & 0\\ \tau (a^{})& \sigma \end{array}\right)$$
lie in the same double coset if and only if $`aa^{}\text{ mod }\pi `$.
###### Proof.
Suppose there exist invertible upper triangular $`2\times 2`$ matrices $`D_1`$, $`D_2`$, $`C_1`$ and $`C_2`$, and $`2\times 2`$ matrices $`X`$ and $`Y`$ such that
$$\left(\begin{array}{cc}D_1& X\\ 0& D_2\end{array}\right)\left(\begin{array}{cc}\sigma & 0\\ \tau (a)& \sigma \end{array}\right)\left(\begin{array}{cc}C_1& Y\\ 0& C_2\end{array}\right)=\left(\begin{array}{cc}\sigma & 0\\ \tau (a^{})& \sigma \end{array}\right),$$
then the following identities must hold:
(8) $`D_1\sigma C_1+X\tau (a)C_1`$ $`=`$ $`\sigma ,`$
(9) $`D_1\sigma Y+X\tau (a)Y+X\sigma C_2`$ $`=`$ $`0,`$
(10) $`D_2\tau (a)C_1`$ $`=`$ $`\tau (a^{}),`$
(11) $`D_2\tau (a)Y+D_2\sigma C_2`$ $`=`$ $`\sigma .`$
Now, reducing (8) modulo $`\pi `$ gives
$$D_1\sigma \sigma C_1^1mod\pi .$$
Comparing the top-left and bottom-right entries of both sides shows that $`D_1`$ and $`C_1`$ are in fact congruent to diagonal matrices modulo $`\pi `$. Similarly, (11) can be used to show that $`D_2`$ and $`C_2`$ are also congruent to diagonal matrices modulo $`\pi `$. Thus $`D_2`$ can be written in the form $`\left(\begin{array}{cc}d_{11}& \pi d_{12}\\ 0& d_{22}\end{array}\right)`$ where $`d_{11}`$ and $`d_{22}`$ are units. Similarly, $`C_1`$ can be written in the form $`\left(\begin{array}{cc}c_{11}& \pi c_{12}\\ 0& c_{22}\end{array}\right)`$ where $`c_{11}`$ and $`c_{22}`$ are units. Substituting in (10) and comparing entries gives
$$d_{11}c_{11}d_{22}c_{22}d_{22}c_{11}1\text{ mod }\pi ,$$
$$d_{11}c_{22}aa^{}\text{ mod }\pi .$$
The above equations are readily seen to imply that $`aa^{}\text{ mod }\pi `$. ∎
###### Corollary 5.2.
When $`n=4`$ and $`k=2`$, there exist at least $`|𝐤|`$ double cosets in $`B\backslash G/B`$ that lie above the permutation $`w`$.
Remark. Note that the intersection invariants do not distinguish between all the double cosets. Among the matrices considered in Proposition 5.1, the double coset corresponding to $`a1\text{ mod }\pi `$ can be distinguished from the others by the $`(2,3)`$th entry of the matrix of intersection numbers described in §2.3. However, the intersection types and intersection numbers do not distinguish between the remaining $`(|𝐤|1)`$ double cosets.
## 6. The general picture
###### Theorem 6.1.
The number of double cosets in $`B\backslash GL_n(A)/B`$ does not depend on $`𝐤`$ if and only if $`n2`$ or $`(n,k)=(3,2)`$ or $`k=1`$.
###### Proof.
We have seen in Sections 3 and 4 that the number of double cosets does not depend on $`𝐤`$ when $`n2`$ or when $`n=3`$ and $`k=2`$. We have also seen in Sections 4 and 5 that for $`n=3`$ and $`k>2`$ and for $`n=4`$ and $`k>1`$ the number of double cosets does depend on $`𝐤`$. In any other case $`(n,k)`$ there exists $`n^{}<n`$ such that the number of cosets depends on $`𝐤`$ for the case $`(n^{},k)`$. Take a permutation matrix $`w^{}`$ of order $`n^{}`$ for which the number of double cosets with permutation invariant $`w^{}`$ in $`GL_n^{}(A)`$ depends on $`𝐤`$. Let $`w`$ be the block matrix
$$\left(\begin{array}{cc}0& I_{nn^{}}\\ w^{}& 0\end{array}\right),$$
where $`I_{nn^{}}`$ is the identity matrix of order $`nn^{}`$. By Section 2.4, $`N(w)N(w^{})`$ so that the number of double cosets with permutation invariant $`w`$ also depends on $`𝐤`$. ∎
The complexity of the double coset space is sketched in the following table
$$\begin{array}{ccccccccccc}& & 1& 2& 3& 4& \mathrm{}& n& & & \\ & & & & & & & & & & \\ 1& & D& D& D& D& \mathrm{}& & & & \\ 2& & D& D& D& N& \mathrm{}& & & & \\ 3& & D& D& N& N& \mathrm{}& & & & \\ 4& & D& D& N& N& \mathrm{}& & & & \\ 5& & D& D& N& N& \mathrm{}& & & & \\ \mathrm{}& & \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& & & & & \\ k& & & & & & & & & & \end{array}$$
Here $`N`$ stands for non-discrete and $`D`$ stands for discrete, to indicate whether the double coset space depends or does not depend (respectively) on the ring in question, e.g., on the cardinality of the residue field.
> Uri Onn
> Einstein Institute of Mathematics, Edmond Safra Campus, Givat Ram,
> Jerusalem 91904, Israel
> urion@math.huji.ac.il
>
> Amritanshu Prasad
> The Institute of Mathematical Sciences, CIT campus,
> Chennai 600 113, India
> amri@imsc.res.in
>
> Leonid Vaserstein
> Department of Mathematics, Penn State University, University Park PA
> 16802-6401, USA
> vstein@math.psu.edu |
warning/0506/nlin0506033.html | ar5iv | text | # Exact Analysis of the Adiabatic Invariants in Time-Dependent Harmonic Oscillator
## Abstract
The theory of adiabatic invariants has a long history and important applications in physics but is rarely rigorous. Here we treat exactly the general time-dependent 1-D harmonic oscillator, $`\ddot{q}+\omega ^2(t)q=0`$ which cannot be solved in general. We follow the time-evolution of an initial ensemble of phase points with sharply defined energy $`E_0`$ and calculate rigorously the distribution of energy $`E_1`$ after time $`T`$, and all its moments, especially its average value $`\overline{E_1}`$ and variance $`\mu ^2`$. Using our exact WKB-theory to all orders we get the exact result for the leading asymptotic behaviour of $`\mu ^2`$.
Adiabatic invariants, usually denoted by $`I`$, in time dependent dynamical systems (not necessarily Hamiltonian), are approximately conserved during a slow process of changing system parameters over a long typical time scale $`T`$. This statement is asymptotic in the sense that the conservation is exact in the limit $`T\mathrm{}`$, whilst for finite $`T`$ we see the deviation $`\mathrm{\Delta }I=I_fI_i`$ of final value of $`I_f`$ from its initial value $`I_i`$ and would like to calculate $`\mathrm{\Delta }I`$. Here we just remind that for the one-dimensional harmonic oscillator it is known since Ehrenfest that the adiabatic invariant for $`T=\mathrm{}`$ is $`I=E/\omega `$, which is the ratio of the total energy $`E=E(t)`$ and the frequency of the oscillator $`\omega (t)`$, both being a function of time. Of course, $`2\pi I`$ is exactly the area in the phase plane $`(q,p)`$ enclosed by the energy contour of constant $`E`$. A general introductory account can be found in Robnik (2005) and references therein, especially Landau and Lifshitz (1996); Reinhardt (1994). However, in the literature this $`\mathrm{\Delta }I`$ is not even precisely defined. As a consequence of that there is considerable confusion about its meaning. Let us just mention the case of periodic parametric resonance in one-dimensional harmonic oscillator where the driving is periodic and yet e.g. the total energy of the system can grow indefinitely for certain system parameter values. (In this work we give a precise meaning to these and similar statements.) Therefore to be on rigorous side we must carefully define what we mean by $`\mathrm{\Delta }I`$. This can be done by considering an ensemble of initial conditions at time $`t=0`$ just before the adiabatic process starts. Of course, there is a vast freedom in choosing such ensembles. In an integrable conservative Hamiltonian system the most natural and the most important choice is taking as the initial ensemble all phase points uniformly distributed on the initial $`N`$-torus, uniform w.r.t. the angle variables. Such an ensemble has a sharply defined initial energy $`E_0`$. Then we let the system evolve in time, not necessarily slowly, and calculate the probability distribution $`P(E_1)`$ of the final energy $`E_1`$, or of other dynamical quantities.
This is in general a difficult problem, but in this work we confine ourself to the one-dimensional general time-dependent harmonic oscillator, so $`N=1`$, described by the Newton equation
$$\ddot{q}+\omega ^2(t)q=0$$
(1)
and work out rigorously $`P(E_1)`$. Given the general dependence of the oscillator’s frequency $`\omega (t)`$ on time $`t`$ the calculation of $`q(t)`$ is already a very difficult unsolvable problem. In the sense of mathematical physics (1) is exactly equivalent to the one-dimensional stationary Schrödinger equation: the coordinate $`q`$ appears instead of the probability amplitude $`\psi `$, time $`t`$ appears instead of the coordinate $`x`$ and $`\omega ^2(t)`$ plays the role of $`EV(x)`$ = energy – potential. In this paper we solve the above stated problem for the general one-dimensional harmonic oscillator, but the details of our calculations are delegated to another publication Robnik and Romanovski (2005).
We begin by defining the system by giving its Hamilton function $`H=H(q,p,t)`$, whose numerical value $`E(t)`$ at time $`t`$ is precisely the total energy of the system at time $`t`$, and for the one-dimensional harmonic oscillator this is
$$H=\frac{p^2}{2M}+\frac{1}{2}M\omega ^2(t)q^2,$$
(2)
where $`q,p,M,\omega `$ are the coordinate, the momentum, the mass and the frequency of the linear oscillator, respectively. The dynamics is linear in $`q,p`$, as described by (1), but nonlinear as a function of $`\omega (t)`$ and therefore is subject to the nonlinear dynamical analysis. By using the index $`0`$ and $`1`$ we denote the initial ($`t=t_0`$) and final $`(t=t_1)`$ value of the variables, and by $`T=t_1t_0`$ we denote the time interval of changing the parametrs of the system.
We consider the phase flow map (we shall call it transition map)
$$\mathrm{\Phi }:\left(\begin{array}{c}q_0\\ p_0\end{array}\right)\left(\begin{array}{c}q_1\\ p_1\end{array}\right).$$
(3)
Because equations of motion are linear in $`q`$ and $`p`$, and since the system is Hamiltonian, $`\mathrm{\Phi }`$ is a linear area preserving map, that is,
$$\mathrm{\Phi }=\left(\begin{array}{cc}a& b\\ c& d\end{array}\right),$$
(4)
with $`det(\mathrm{\Phi })=adbc=1`$. Let $`E_0=H(q_0,p_0,t=t_0)`$ be the initial energy and $`E_1=H(q_1,p_1,t=t_1)`$ be the final energy, that is,
$$\begin{array}{c}E_1=\frac{1}{2}\left(\frac{(cq_0+dp_0)^2}{M}+M\omega _1^2(aq_0+bp_0)^2\right).\hfill \end{array}$$
(5)
Introducing the new coordinates, namely the action $`I=E/\omega `$ and the angle $`\varphi `$,
$$q_0=\sqrt{\frac{2E_0}{M\omega _0^2}}\mathrm{cos}\varphi ,p_0=\sqrt{2ME_0}\mathrm{sin}\varphi $$
(6)
from (5) we obtain
$$E_1=E_0(\alpha \mathrm{cos}^2\varphi +\beta \mathrm{sin}^2\varphi +\gamma \mathrm{sin}2\varphi ),$$
(7)
where
$$\begin{array}{c}\alpha =\frac{c^2}{M^2w_0^2}+a^2\frac{\omega _1^2}{\omega _0^2},\beta =d^2+\omega _1^2M^2b^2,\hfill \\ \hfill \gamma =\frac{cd}{M\omega _0}+abM\frac{\omega _1^2}{\omega _0}.\end{array}$$
(8)
Given the uniform probability distribution of initial angles $`\varphi `$ equal to $`1/(2\pi )`$, which defines our initial ensemble at time $`t=0`$, we can now calculate the averages. Thus
$$\overline{E}_1=\frac{1}{2\pi }E_1𝑑\varphi =\frac{E_0}{2}(\alpha +\beta ).$$
(9)
That yields $`E_1\overline{E}_1=E_0(\delta \mathrm{cos}2\varphi +\gamma \mathrm{sin}2\varphi )`$ and
$$\mu ^2=\overline{(E_1\overline{E}_1)^2}=\frac{E_0^2}{2}\left(\delta ^2+\gamma ^2\right),$$
(10)
where we have denoted $`\delta =(\alpha \beta )/2`$.
It follows from (8), (9) that we can write (10) also in the form
$$\mu ^2=\overline{(E_1\overline{E}_1)^2}=\frac{E_0^2}{2}\left[\left(\frac{\overline{E}_1}{E_0}\right)^2\frac{\omega _1^2}{\omega _0^2}\right].$$
(11)
It is straightforward to show that for arbitrary positive integer $`m`$, we have $`\overline{(E_1\overline{E}_1)^{2m1}}=0`$ and
$$\overline{(E_1\overline{E}_1)^{2m}}=\frac{(2m1)!!}{m!}\left(\overline{(E_1\overline{E}_1)^2}\right)^m.$$
(12)
Thus $`2m`$-th moment of $`P(E_1)`$ is equal to $`(2m1)!!\mu ^{2m}/m!`$, and therefore, indeed, all moments of $`P(E_1)`$ are uniquely determined by the first moment $`\overline{E_1}`$.
Expression (11) is positive definite by definition and this leads to the first interesting conclusion: In full generality (no restrictions on the function $`\omega (t)`$!) we have always $`\overline{E_1}E_0\omega _1/\omega _0`$ and therefore the final value of the adiabatic invariant $`I_1=\overline{E_1}/\omega _1`$ is always greater or equal to the initial value $`I_0=E_0/\omega _0`$. In other words, the value of the adiabatic invariant never decreases, which is a kind of irreversibility statement. Moreover, it is constant only for infinitely slow processes $`T=\mathrm{}`$, which is an ideal adiabatic process, i.e. $`\mu =0`$. For periodic processes $`\omega _1=\omega _0`$ we see that always $`\overline{E_1}E_0`$, so the mean energy never decreases. The other extreme to $`T=\mathrm{}`$ is the instantaneous ($`T=0`$) jump where $`\omega _0`$ switches to $`\omega _1`$ discontinuously, whilst $`q`$ and $`p`$ remain continuous, and this results in $`a=d=1`$ and $`b=c=0`$, and then we find
$$\overline{E_1}=\frac{E_0}{2}(\frac{\omega _1^2}{\omega _0^2}+1),\mu ^2=\frac{E_0^2}{8}\left[\frac{\omega _1^2}{\omega _0^2}1\right]^2.$$
(13)
Below we shall treat the special case with $`\omega _1^2=2\omega _0^2`$, and thus will find $`\mu ^2/E_0^2=1/8=0.125`$.
Our general study now focuses on the calculation of the transition map (4), namely its matrix elements $`a,b,c,d`$. Starting from the Hamilton function (2) and its Newton equation (1) we consider two linearly independent solutions $`\psi _1(t)`$ and $`\psi _2(t)`$ and introduce the matrix
$$\mathrm{\Psi }(t)=\left(\begin{array}{cc}\psi _1(t)& \psi _2(t)\\ M\dot{\psi }_1(t)& M\dot{\psi }_2(t)\end{array}\right).$$
(14)
Consider a solution $`\widehat{q}(t)`$ of (1) such that
$$\widehat{q}(t_0)=q_0,\dot{\widehat{q}}(t_0)=p_0/M.$$
(15)
Because $`\psi _1`$ and $`\psi _2`$ are linearly independent, we can look for $`\widehat{q}(t)`$ in the form
$$\widehat{q}(t)=A\psi _1(t)+B\psi _2(t).$$
(16)
Then $`A`$ and $`B`$ are determined by
$$\left(\begin{array}{c}A\\ B\end{array}\right)=\mathrm{\Psi }^1(t_0)\left(\begin{array}{c}q_0\\ p_0\end{array}\right).$$
(17)
Let $`q_1=\widehat{q}(t_1),p_1=M\dot{\widehat{q}}(t_1)`$. Then from (15)–(17) we see that
$$\left(\begin{array}{c}q_1\\ p_1\end{array}\right)=\mathrm{\Psi }(t_1)\mathrm{\Psi }^1(t_0)\left(\begin{array}{c}q_0\\ p_0\end{array}\right).$$
(18)
We recognize the matrix on the right-hand side of (18) as the transition map $`\mathrm{\Phi }`$, that is,
$$\mathrm{\Phi }=\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)=\mathrm{\Psi }(t_1)\mathrm{\Psi }^1(t_0).$$
(19)
Due to lack of space we mention only one specific model, namely the linear model defined by the piecewise linear $`\omega ^2(t)`$ function
$$\omega ^2(t)=\{\begin{array}{c}\omega _0^2\mathrm{if}t0\hfill \\ \omega _0^2+\frac{(\omega _1^2\omega _0^2)}{T}t\mathrm{if}0<t<T\hfill \\ \omega _1^2\mathrm{if}tT\hfill \end{array}.$$
(20)
In this case the equation (1) can be solved exactly in terms of the Airy functions, and the formalism explained above leads in a straightforward but lengthy manner to the final exact result for $`\overline{E_1}`$ and consequently for $`\mu ^2`$ etc. It is too complex to be shown here. The special case $`\omega _0^2=1`$ and $`\omega _1^2=2`$ has been checked very carefully, also numerically, and $`\mu ^2`$ goes correctly from $`1/8`$ at $`T=0`$ to zero as $`T\mathrm{}`$, in a typical oscillatory way. Using the well known asymptotic expressions for the Airy functions we find the leading asymptotic approximation
$$\frac{\mu ^2}{E_0^2}=\frac{\overline{(E_1\overline{E}_1)^2}}{E_0^2}\frac{ϵ^2}{128}\left(94\sqrt{2}\mathrm{cos}(\frac{48\sqrt{2}}{3ϵ})\right),$$
(21)
where we introduce the adiabatic parameter $`ϵ=1/T`$ which is assumed small. Please observe the oscillatory approach to zero as $`ϵ0`$, which in the mean goes to zero quadratically as $`ϵ^2`$.
Returning to the general case we now mention that the final energy distribution function written down as
$$P(E_1)=\frac{1}{2\pi }\underset{j=1}{\overset{4}{}}\left|\frac{d\varphi }{dE_1}\right|_{\varphi =\varphi _j(E_1)}$$
(22)
cannot be calculated analytically in a closed form in any useful way, because it boils down to finding the roots of a quartic polynomial, so we do not try to do that here, although numerically it shows interesting aspects. It has a finite interval as its support, between the lower limit $`E_{min}`$ and the upper limit $`E_{max}`$, and at both values it has an integrable singularity of the type $`1/\sqrt{x}`$. In between for every value of $`E_1=const=E_1(\varphi )`$, this equation has four solutions, namely $`\varphi _1,\varphi _2,\varphi _3,\varphi _4`$, and thus we have to sum up all four contributions in the general formula (22). On the other hand, as we have seen, the moments of this interesting distribution function can be calculated exactly to all orders.
We proceed with the calculation of the transition map $`\mathrm{\Phi }`$ in the general case, and because (1) is generally not solvable, we have ultimately to resort to some approximations. Since the adiabatic limit $`ϵ0`$ is the asymptotic regime that we would like to understand, the application of the rigorous WKB theory (up to all orders) is most convenient, and usually it turns out that the leading asymptotic terms are well described by just the leading WKB terms.
We introduce re-scaled and dimensionless time $`\lambda `$
$$\lambda =ϵt,ϵ=1/T$$
(23)
so that (1) is transformed to the equation
$$ϵ^2q^{\prime \prime }(\lambda )+\omega ^2(\lambda )q(\lambda )=0.$$
(24)
Let $`q_+(\lambda )`$ and $`q_{}(\lambda )`$ be two linearly independent solutions of (24). Then the matrix (14) takes the form
$$\mathrm{\Psi }_\lambda =\left(\begin{array}{cc}q_+(\lambda )& q_{}(\lambda )\\ ϵMq_+^{}(\lambda )& ϵMq_{}^{}(\lambda )\end{array}\right)$$
(25)
and taking into account that $`\lambda _0=ϵt_0,\lambda _1=ϵt_1`$ we obtain for the matrix (19) the expression
$$\mathrm{\Phi }=\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)=\mathrm{\Psi }_\lambda (\lambda _1)\mathrm{\Psi }_\lambda ^1(\lambda _0).$$
(26)
We now use the WKB method in order to obtain the coefficients $`a,b,c,d`$ of the matrix $`\mathrm{\Phi }`$. To do so, we look for solution of (24) in the form
$$q(\lambda )=w\mathrm{exp}\left\{\frac{1}{ϵ}\sigma (\lambda )\right\}$$
(27)
where $`\sigma (\lambda )`$ is a complex function that satisfies the differential equation
$$(\sigma ^{}(\lambda ))^2+ϵ\sigma ^{\prime \prime }(\lambda )=\omega ^2(\lambda )$$
(28)
and $`w`$ is some constant with dimension of length. The WKB expansion for the phase is
$$\sigma (\lambda )=\underset{k=0}{\overset{\mathrm{}}{}}ϵ^k\sigma _k(\lambda ).$$
(29)
Substituting (29) into (28) and comparing like powers of $`ϵ`$ gives the recursion relation
$$\sigma _0^2=\omega ^2(\lambda ),\sigma _n^{}=\frac{1}{2\sigma _0^{}}(\underset{k=1}{\overset{n1}{}}\sigma _k^{}\sigma _{nk}^{}+\sigma _{n1}^{\prime \prime }).$$
(30)
Here we apply our WKB notation and formalism Robnik and Romanovski (2000) and we can choose $`\sigma _{0,+}^{}(\lambda )=\mathrm{i}\omega (\lambda )\mathrm{or}\sigma _{0,}^{}(\lambda )=\mathrm{i}\omega (\lambda )`$. That results in two linearly independent solutions of (24) given by the WKB expansions with the coefficients
$$\begin{array}{c}\sigma _{0,\pm }(\lambda )=\pm \mathrm{i}_{\lambda _0}^\lambda \omega (x)𝑑x,\sigma _{1,\pm }(\lambda )=\frac{1}{2}\mathrm{log}\frac{\omega (\lambda )}{\omega (\lambda _0)},\hfill \\ \hfill \sigma _{2,\pm }=\pm \frac{\mathrm{i}}{8}_{\lambda _0}^\lambda \frac{3\omega ^{}(x)^22\omega (x)\omega ^{\prime \prime }(x)}{\omega (x)^3}𝑑x,\mathrm{}\end{array}$$
(31)
Since $`\omega (\lambda )`$ is a real function we deduce from (30) that all functions $`\sigma _{2k+1}^{}`$ are real and all functions $`\sigma _{2k}^{}`$ are pure imaginary and $`\sigma _{2k,+}^{}=\sigma _{2k,}^{},\sigma _{2k+1,+}^{}=\sigma _{2k+1,}^{}`$ where $`k=0,1,2,\mathrm{}`$, and thus we have $`\sigma _+^{}=A(\lambda )+\mathrm{i}B(\lambda ),\sigma _{}^{}=A(\lambda )\mathrm{i}B(\lambda )`$ where $`A(\lambda )=_{k=0}^{\mathrm{}}ϵ^{2k+1}\sigma _{2k+1}^{}(\lambda ),B(\lambda )=\mathrm{i}_{k=0}^{\mathrm{}}ϵ^{2k}\sigma _{2k,+}^{}(\lambda )`$. Integration of the above equations yields
$$\sigma _+=r(\lambda )+\mathrm{i}s(\lambda ),\sigma _{}=r(\lambda )\mathrm{i}s(\lambda ),$$
(32)
where $`r(\lambda )=_{\lambda _0}^\lambda A(x)𝑑x,s(\lambda )=_{\lambda _0}^\lambda B(x)𝑑x`$. Below we shall denote $`s_1=s(\lambda _1)`$.
Using this notation we find that the elements of the transition matrix $`\mathrm{\Phi }`$ have the following form, after taking into account that $`det(\mathrm{\Phi })=abcd=1`$,
$`a=`$ $`{\displaystyle \frac{1}{\sqrt{B_0B_1}}}\left[A_0\mathrm{sin}\left({\displaystyle \frac{s_1}{ϵ}}\right)B_0\mathrm{cos}\left({\displaystyle \frac{s_1}{ϵ}}\right)\right],`$ (33)
$`b=`$ $`{\displaystyle \frac{1}{M\sqrt{B_0B_1}}}\mathrm{sin}\left({\displaystyle \frac{s_1}{ϵ}}\right),`$
$`c=`$ $`{\displaystyle \frac{M}{\sqrt{B_0B_1}}}[(A_0A_1+B_0B_1)\mathrm{sin}\left({\displaystyle \frac{s_1}{ϵ}}\right)`$
$`+(A_0B_1A_1B_0)\mathrm{cos}\left({\displaystyle \frac{s_1}{ϵ}}\right)],`$
$`d=`$ $`{\displaystyle \frac{1}{\sqrt{B_0B_1}}}\left[A_1\mathrm{sin}\left({\displaystyle \frac{s_1}{ϵ}}\right)+B_1\mathrm{cos}\left({\displaystyle \frac{s_1}{ϵ}}\right)\right].`$
This is so far exact result, based on the WKB expansion technique. What we are mostly interested in is the asymptotic behaviour of $`\mu ^2`$ when $`ϵ`$ is small and tends to zero. All other aspects and technical details will be published in a separate paper Robnik and Romanovski (2005).
Let us consider the first order WKB approximation, that is,
$$A(\lambda )ϵ\sigma _{1,+}^{}(\lambda ),B(\lambda )\frac{\sigma _{0,+}^{}(\lambda )}{\mathrm{i}}=\omega (\lambda ).$$
(34)
We find for the variance (11)
$$\begin{array}{c}\frac{\mu ^2}{E_0^2}=ϵ^2(\frac{\omega _1^2\omega _{0}^{}{}_{}{}^{2}}{8\omega _0^6}+\frac{\omega _{1}^{}{}_{}{}^{2}}{8\omega _0^2\omega _1^2}\hfill \\ \hfill \frac{\omega _0^{}\omega _1^{}}{4\omega _0^4}\mathrm{cos}\left(\frac{2}{ϵ}_{\lambda _0}^{\lambda _1}\omega (x)dx\right))+O(ϵ^3).\end{array}$$
(35)
Substituting into the last formula $`\omega (\lambda )=\sqrt{1+\lambda }`$ we obtain exactly the approximation (21).
Suppose now that all derivatives at $`\lambda _0`$ and $`\lambda _1`$ vanish up to order $`(n1)`$, i.e. $`\omega ^{}(\lambda _0)=\omega ^{}(\lambda _1)=\mathrm{}=\omega ^{(n1)}(\lambda _0)=\omega ^{(n1)}(\lambda _1)=0`$, and $`\omega ^{(n)}(\lambda _0)\omega ^{(n)}(\lambda _1)0`$. Then $`\sigma _1^{}(\lambda _0)=\sigma _1^{}(\lambda _1)=\mathrm{}=\sigma _{n1}^{}(\lambda _0)=\sigma _{n1}^{}(\lambda _1)=0,\sigma _n^{}(\lambda _0)\sigma _n(\lambda _1)0.`$
Hence, in the case $`n=2k1`$ we can assume
$`A(\lambda )=`$ $`ϵ^{2k1}\sigma _{2k1,+}^{}(\lambda )+h.o.t.`$ (36)
$`B(\lambda )=`$ $`\omega (\lambda )\mathrm{i}ϵ^{2k}\sigma _{2k,+}^{}(\lambda )+h.o.t.`$
and obtain
$$\begin{array}{c}\frac{\mu ^2}{E_0^2}=ϵ^{4k2}(\frac{\sigma _{2k1,+}^{}(\lambda _1)^2}{2\omega _0^2}+\frac{\omega _1^2\sigma _{2k1,+}^{}(\lambda _0)^2}{2\omega _0^4}\hfill \\ \hfill \frac{\omega _1\sigma _{2k1,+}^{}(\lambda _0)\sigma _{2k1,+}^{}(\lambda _1)}{\omega _0^3}\mathrm{cos}\left(\frac{2s_1}{ϵ}\right))\\ \hfill +O(ϵ^{4k1}).\end{array}$$
(37)
In the case when $`n=2k`$ we can suppose
$`A(\lambda )=`$ $`ϵ^{2k+1}\sigma _{2k+1,+}^{}(\lambda )+h.o.t.`$ (38)
$`B(\lambda )=`$ $`\omega (\lambda )\mathrm{i}ϵ^{2k}\sigma _{2k,+}^{}(\lambda )+h.o.t.`$
Then, similarly as above, we obtain
$$\begin{array}{c}\frac{\mu ^2}{E_0^2}=ϵ^{4k}(\frac{\sigma _{2k,+}^{}(\lambda _1)^2}{2\omega _0^2}+\frac{\omega _1^2\sigma _{2k,+}^{}(\lambda _0)^2}{2\omega _0^4}\hfill \\ \hfill \frac{\omega _1\sigma _{2k,+}^{}(\lambda _0)\sigma _{2k,+}^{}(\lambda _1)}{\omega _0^3}\mathrm{cos}\left(\frac{2s_1}{ϵ}\right))+O(ϵ^{4k+1}).\end{array}$$
(39)
From this we can conclude that if $`\omega (t)`$ is of class $`𝒞^m`$ (having $`m`$ continuous derivatives, $`m=n1`$) $`\mu ^2`$ goes to zero oscillating but in the mean as $`ϵ^{2n}=ϵ^{2(m+1)}`$. If $`m=\mathrm{}`$ (analytic functions) according to Landau and Lifshitz Landau and Lifshitz (1996) the decay to zero is oscillating and on the average is exponential $`\mathrm{exp}(const/ϵ)`$.
###### Acknowledgements.
This work was supported by the Ministry of Higher Education, Science and Technology of the Republic of Slovenia, Nova Kreditna Banka Maribor and Telekom Slovenije. |
warning/0506/cond-mat0506680.html | ar5iv | text | # Universal properties of three-dimensional magnetohydrodynamic turbulence: Do Alfvén waves matter?
## I Introduction
The effects of propagating waves on the statistical properties of systems out of equilibrium remain an important topic of discussion. In the context of a coupled spin model in one-dimension $`(d=1)`$dib it has been shown that the presence of such waves leads to weak dynamic scaling in that model. In contrast, in a coupled Burgers-like model in $`d=1`$ propagating waves do not affect the scaling properties of the correlation functions at all abjkb . So far such issues have been considered only within very simplified one-dimensional nonequilibrium models abjkb ; verma . Magnetohydrodynamic ($`3d`$MHD) turbulence, which is a hydrodynamic description of the coupled evolution of the velocity fields $`𝐮`$ and the magnetic fields $`𝐛`$ in a quasi-neutral plasma, stands as a very good candidate for a natural system with propagating waves in three dimensions as most of its natural realizations have proagating Alfvén waves arising due to the presence of mean magnetic fields. Examples of such physical situations include solar wind, neutral plasma in fusion confinement devices etc. The presence of the propagating Alfvén modes, in addition to the usual dissipative modes due to the fluid and magnetic viscosities makes MHD turbulence a good natural example to study the interplay between the propagating and the dissipative modes in a system and their combined effects on the scaling properties of the correlation and structure functions, which are important issues from the point of view of nonequilibrium statistical mechanics.
The scaling of magnetic- \[$`E_b(k)`$\] and kinetic- \[$`E_u(k)`$\] energy spectra in the inertial range (i.e., wavevector $`k`$ lies in the region $`L^1k\eta _D^1`$, $`L`$ and $`\eta _D`$ being the integral scale given by the system size and the dissipation scale, respectively) in $`3d`$MHD in the presence of Alfvén waves originating due to a mean magnetic field $`𝐁_𝐨`$ remains controversial till the date. Numerical simulations, due to the lack of sufficient resolutions failed to conclusively distinguish between the Kolmogorov’s and Kraichnan’s predictions (see below). In this communication, within one-loop self-consistent mode-coupling (SCMC) approximations, we obtain the following results for homogeneous but anisotropic (due to the mean magnetic field) $`3d`$MHD:
* The bare Alfvén wave speed, proportional to $`B_o`$, renormalizes to acquire a singular $`k`$-dependence to become $`B_ok^{1/3}`$ where $`k`$ is a Fourier wavevector belonging to the inertial (scaling) range. From this result we are able to conclude that even in the presence of a mean magnetic field the kinetic- and the magnetic- energy spectra scale as $`k^{5/3}`$ in the inertial range, identical with the situation without any mean magnetic field.
* The dimensionless Kolmogorov’s constants for the kinetic- and magnetic-energy spectra depends on a dimensionless parameter $`\beta `$ which we identify as the ratio of the renormalised Alfvén wave speed and the renormalised viscosities (see below).
* The intermittency exponents of the Elsässer fields $`𝐳^\pm =𝐮\pm 𝐛`$ which approximately and qualitatively characterise the multiscaling properties of the structure functions in a log-normal model depend on the parameter $`\beta `$ and decrease with increasing $`\beta `$.
Thus our results show that Alfvén waves in $`3d`$MHD do not affect the scaling properties of the two-point correlation functions. However, the multiscaling properties of the structure functions are shown to be affected by the mean magnetic fields.
The famous Kolmogorov’s arguments k41 for fluid turbulence can be easily extended to $`3d`$MHD turbulence. In MHD, in the unit where mass density $`\rho =1`$, the kinetic energy dissipation rate per unit mass ($`ϵ_K`$) and the magnetic energy dissipation rate per unit mass ($`ϵ_M`$) have same physical dimensions. Thus, as in fluid turbulence, by claiming that the structure functions of the velocity and magnetic field differences in the inertial range must be constructed out of the mean energy dissipation rate $`ϵ`$ per unit mass and the local length scale $`r`$ (belonging to the inertial range) one obtains for the $`n`$-th order structure function k41mhd ; abthesis
$$S_n^a(r)[a_i(𝐱+𝐫)a_i(𝐱)]\widehat{r}_i]^n(ϵr)^{n/3},\eta _DrL$$
(1)
where $`ϵ=ϵ_K\mathrm{or}ϵ_M`$ and $`a=u,b`$ for the velocity and the magnetic fields, $`\eta _D`$ is the (small) dissipation scale and $`L`$ is the (large) system size. This yields, for the energy spectra in the inertial range (as a function of wavevector $`k`$)
$$E_a(k)k^{5/3},a=u,b.$$
(2)
This is known as the K41 theory in the relevant literature. However, in the presence of a mean magnetic field $`𝐁_𝐨`$, the Alfvén waves are generated (see below) with the propagation speed $`B_o`$. Thus in such a case in addition to the usual Kolmogorov time scale $`k^{2/3}`$ yakhot ; abepl , there exists another time scale constructed out of the mean magnetic field $`B_o`$, known as the Kraichnan time scale $`(B_ok)^1`$ krai . Kraichnan argued that this time scale would determine the energy cascade process and hence would enter in the expression of the structure function yielding $`E_a(k)k^{3/2},a=u,b`$. There has not been any satisfactory resolution of this issue till this date; due to the particularly difficult vectorial nature of the $`3d`$MHD equations (see below) it is rather difficult to achieve high Reynolds number in Direct Numerical Solutions (DNS) of the $`3d`$MHD equations. Numerical solutions of an MHD shell-model in the presence of a small mean magnetic field-like term did not find any dependence of the multiscaling of the structure functions on the mean magnetic field abthesis . Analytically, it has been shown within the context of a $`1d`$ coupled Burgers-like model abjkb that the energy spectra are independent of a mean magnetic field and in case of $`1d`$ MHD turbulence it exhibits the Kolmogorov scaling for the energy spectra. Similar conclusions followed from an analogous one-dimensional model verma . Various phenomenological approaches, including weak turbulence theories and three-wave interaction models, yield, in general, mean magnetic field dependences of the energy spectra galtier ; ng ; montgomery ; shebalin when the mean magnetic field is strong. However, these theories, despite their predictions are either not directly derivable from the underlying $`3d`$MHD Eqs. of motion or involve additional assumptions on the flow fields. Analyses starting from the $`3d`$MHD Eqs. in this regard are still lacking. Most observational results on astrophysical systems seem to favour the K41 results obser . Simulations of incompressible MHD simu1 find results close to the K41 result. Simulations of compressible MHD simu2 , even though in some cases find energy spectra closer to K41, in general yields a less clear picture. Recent numerical results of Müller et al biskampnew suggest that in presence of finite magnetic helicity structure functions parallel and perpendicular to the mean magnetic fields are affected differently by the mean magnetic field. In this paper, we address some of these issues by starting from the $`3d`$MHD Eqs. without making any further assumptions on the velocity and the magnetic fields, except for the validity of the perturbative approaches. In particular, we show by applying one-loop mode coupling methods on the $`3d`$MHD equations with a mean magnetic field and in the absence of any magnetic helicity that the one-dimensional energy spectra in $`3d`$MHD turbulence are independent of the mean magnetic field $`𝐁_𝐨`$ and scale as $`k^{5/3}`$ in the inertial range where $`𝐤`$ is a Fourier wavevector belonging to the inertial range. We then proceed to calculate the Kolmogorov’s constants for the kinetic- and the magnetic-energy spectra and show that they depend on $`B_o`$. Lastly, we calculate the intermittency exponents for the velocity and the magnetic fields and find that they decrease with increasing $`B_o`$. We do not distinguish between the longitudinal and the transverse structure functions. Even though for analytical convenience we assume a weak mean magnetic field, we are able to obtain new and interesting results concerning scaling and multiscaling in $`3d`$MHD in the presence of a mean magnetic field of small magnitude. In this respect our results can be considered as complementary to some of the existing results galtier ; montgomery ; ng . The rest of the paper is organised as follows: In Sec.II we discuss the $`3d`$MHD equations for incompressible fluids. In Sec.V we show that for incompressible fluids the bare Alfvén wave speed $`B_o`$ renormalizes to pick a correction $`k^{1/3}`$ in the inertial range. This, as we argue, implies that the energy spectra scale as $`k^{5/3}`$ in the inertial range even in the presence of a mean magnetic field. In Sec.VI we calculate the Kolmogorov’s constants for the kinetic and the magnetic energy spectra. We introduce a parameter $`\beta `$ which is the dimensionless ratio of the renormalised mean magnetic field and renormalised viscosities and show that the Kolmogorov’s constants depend on $`\beta `$. In Sec.VII we elucidate the multiscaling properties of the structure functions by calculating the intermittency exponents and show that they decrease with increasing $`\beta `$, i.e., with increasing $`B_o`$. This suggests that a mean magnetic field tends to reduce multiscaling corrections to the K41 results for the structure functions. In Sec.VIII we summarize our results.
## II Model equations
We begin by writing down the $`3d`$MHD equations for the velocity fields $`𝐮`$ and the magnetic fields $`𝐛`$: The velocity field $`𝐮`$ is governed by the Navier-Stokes equation modified by the inclusion of the Lorentz force jack
$$\frac{𝐮}{t}+(𝐮)𝐮=\frac{p}{\rho }+\frac{(\times 𝐛)\times 𝐛}{\rho }+\nu ^2𝐮+𝐟,$$
(3)
and the dynamics of the magnetic field $`𝐛`$ is governed by the Induction equation jack constructed out of the Ohm’s law for a moving frame and the Ampere’s law:
$$\frac{𝐛}{t}=\times (𝐮\times 𝐛)+\eta ^2𝐛+𝐠.$$
(4)
Here, $`\rho `$ is the mass density, $`p`$ is the pressure, $`\nu `$ is the fluid viscosity and $`\eta `$ is the magnetic viscosity (inversely proportional to the electrical conductivity of the fluid medium concerned). Functions $`𝐟`$ and $`𝐠`$ are external forces needed to maintain a statistical steady state. In the present approach these are taken to be stochastic forces. We assume them to be zero-mean and Gaussian distributed with specified variances (see below). In addition to Eqs. (3) and (4) we also have $`𝐛=\mathrm{𝟎}`$ (Maxwell’s equation) and, for incompressible fluids $`𝐮=\mathrm{𝟎}`$.
If the magnetic fields $`𝐛(𝐱,t)`$ are such that $`𝐛=𝐁_𝐨`$ (a constant vector) then replacing $`𝐛`$ by $`𝐛+𝐁_𝐨`$ where now $`𝐛=\mathrm{𝟎}`$, in Eqs. (3) and (4) one obtains additional linear terms proportional to wavevector $`𝐤`$ leading to wave-like excitations, known as Alfvén waves. The resulting equations are
$$\frac{𝐮}{t}+\lambda _1(𝐮)𝐮)=\frac{p}{\rho }+\lambda _2\frac{(\times 𝐛)\times 𝐛}{\rho }+\frac{(\times 𝐛)\times 𝐁_𝐨}{\rho }+\nu ^2𝐮+𝐟,$$
(5)
and
$$\frac{𝐛}{t}=\lambda _3\times (𝐮\times 𝐛)+\times (𝐮\times 𝐁_𝐨)+\eta ^2𝐛+𝐠.$$
(6)
with $`𝐛(𝐱,𝐭)=\mathrm{𝟎}`$. The parameters $`\lambda _1,\lambda _2,\lambda _3`$ are kept for book keeping purposes and can be set to unity (see Sec.IV). Note that on dropping the nonlinear and the dissipative terms from Eqs. (5) and (6) the resulting linear coupled partial differential equations admit wave-like solutions with dispersion relation linear in wavevector $`k`$. These are known as the Alfvén waves mont in the literature which propagate with speed proportional to $`B_o`$. In the following sections we would calculate the kinetic and the magnetic energy spectra, the Kolmogorov’s constants and the intermittency exponents in the presence of the Alfvén waves, i.e., for $`B_o0`$.
## III Correlation and structure functions in MHD
In the statistical steady state the time dependent correlation functions of $`𝐮`$ and $`𝐛`$ exhibit scaling which are characterized by the roughness exponents $`\chi _u`$ and $`\chi _b`$ respectively, and the dynamic exponent $`z`$. In terms of the scaling exponents $`\chi _u,\chi _b`$ and $`z`$ the velocity and the magnetic field correlators have the form (as a function of wavevector $`𝐤`$ and frequency $`\omega `$)
$`C_{ij}^u(k,\omega )=u_i(𝐤,\omega )u_j(𝐤,\omega )=D^uP_{ij}k^{d2\chi _uz}f_u(k^z/\omega ),`$
$`C_{ij}^b(k,\omega )=b_i(𝐤,\omega )b_j(𝐤,\omega )=D^bP_{ij}k^{d2\chi _bz}f_b(k^z/\omega ),`$ (7)
where $`f_u`$ and $`f_b`$ are scaling functions, $`P_{ij}`$ is the transverse projection operator: $`P_{ij}=\delta _{ij}k_ik_j/k^2`$ to account for the incompressibility of the fields. The kinetic- and magnetic-energy are simply related to the correlation functions: They are just the equal time velocity and magnetic field auto correlation functions multiplied by appropriate phasefactors. We now set out to find whether the Alfvén waves are relevant perturbations on the system. It should be noted that the correlators (7) are chosen as they would be in the fully isotropic case. In the presence of a mean magnetic field $`𝐁_𝐨`$ there would however be additional anisotropic terms in the expression for the correlators above. In the small $`𝐁_𝐨`$ limit the lowest order corrections to the expressions (7) are $`O(B_o)^2`$. Neglecting these correction terms does not affect our scaling analyses below, and in any case, we are interested to see whether anisotropic Alfvén waves are relevant perturbations in the large scale, long time limit on the isotropic $`3d`$MHD as characterised by expressions (7). Therefore, it suffices for us to work with the expressions (7) for our mode-coupling analyses below.
A complete characterisation of the nonequilibrium steady state (NESS) of MHD, however, requires informations about, not just the scaling of the energy-spectra, but also of the $`n`$-th order equal time structure functions of the velocity and the magnetic fields in the inertial range. These are defined by
$$S_n^a(r)|[a_i(𝐱+𝐫)a_i(𝐱)]\widehat{r}_i|^n,a=u,b,\eta _DrL,$$
(8)
which scale as $`r^{\zeta _n^a}`$ where $`r`$ belongs to the inertial range ($`\eta _DrL`$). According to the Kolmogorov theory (K41) the multiscaling exponents $`\zeta _n^a=n/3`$; i.e., they are linear in $`n`$. Subsequent numerical and observational results for homogeneous and isotropic $`3d`$MHD, i.e., without any mean magnetic fields, abprl ; biskamp3d suggested deviations from the K41 results which are similar to those of pure fluid turbulence rahulrev . Much less results for the multiscaling of the structure functions are available when there are mean magnetic fields. Recent numerical simulations of Müller et al biskampnew suggested that structure functions parallel and perpendicular to the mean magnetic fields are differently affected by it when there is a finite magnetic helicity. Below we investigate the issue of multiscaling in presence of a mean magnetic field $`𝐁_𝐨`$ within the context of a log-normal model by evaluating the relevant intermittency exponents which are found to depend explicitly on $`B_o`$. In our analyses we ignore the distinction between the structure functions parallel and perpendicular to the direction of the mean magnetic field; nevertheless, as we discuss below, our results are significantly new.
## IV Symmetries of the equations of motion
We begin by re-expressing $`3d`$MHD Eqs. (3) and (4) by writing the magnetic fields $`𝐛`$ as a sum of a space-time dependent part and a constant vector: $`𝐛(𝐱,t)𝐛(𝐱,t)+\stackrel{~}{𝐁}_𝐨`$. In terms of these fields and parameters, the $`3d`$MHD Eqs. become
$$\frac{𝐮}{t}+\lambda _1(𝐮)𝐮)=\frac{p}{\rho }+\lambda _2\frac{(\times 𝐛)\times 𝐛}{\rho }+\frac{(\times 𝐛)\times \stackrel{~}{𝐁}_𝐨}{\rho }+\nu ^2𝐮+𝐟,$$
(9)
and
$$\frac{𝐛}{t}=\lambda _3\times (𝐮\times 𝐛)+\times (𝐮\times \stackrel{~}{𝐁}_𝐨)+\eta ^2𝐛+𝐠.$$
(10)
It should be noted that in Eqs. (9) and (10) $`\stackrel{~}{𝐁}_𝐨`$ is not the mean magnetic field in general; it would be so only if $`𝐛(𝐱,t)`$ is zero in Eqs. (9) and (10). Note that our above way of splitting the magnetic fields fields does not change the actual mean magnetic field in the system: It is still given by $`\stackrel{~}{𝐁}_𝐨+𝐛=𝐁_𝐨`$. The Eqs. of motion (9) and (10) are invariant under the following continuous transformations abepl ; abpassive :
* The Galilean transformation (TI): $`𝐮(𝐱,t)𝐮(𝐱+𝐮_𝐨t,t)+𝐮_\mathrm{𝟎},\frac{}{t}𝐮_\mathrm{𝟎}.,`$ and $`𝐛𝐛`$ abjkb ; abepl ; fns with $`\lambda _1=\lambda _3=1`$ in Eqs. (5) and (6). This implies non-renormalization of $`\lambda _1`$ abjkb ; fns ; freykpz .
* The transformation (TII) $`\stackrel{~}{𝐁}_\mathrm{𝟎}\stackrel{~}{𝐁}_\mathrm{𝟎}+\lambda _2\delta `$, $`𝐛(𝐱,t)𝐛(𝐱,t)\delta `$, $`𝐮𝐮`$. This allows one to work with the effective magnetic fields defined by $`\sqrt{\lambda }_2𝐛`$ such that the coefficient of the Lorentz force vertex constructed out of the effective magnetic fields does not renormalize. This, therefore, ensures $`\lambda _2`$ can be set to unity abepl ; abpassive by treating all magnetic fields as effective fields. Here the shift $`\delta `$ is a vector. A transformation similar to TII above exists in a problem of passive scalar turbulence anton .
The transformation TII essentially signifies the freedom to split the total magnetic fields as a sum of a constant part and a space time dependent part. It should be noted that the transformation TII keeps the mean magnetic field unchanged. In fact, nonrenormalization of $`\lambda _2`$ can be shown in a simpler way: Let us assume that under mode eliminations and rescaling $`\lambda _2\alpha \lambda _2`$. This scale factor of $`\alpha `$ can now be absorbed by redefining the units of the magnetic fields by $`\sqrt{\alpha }𝐛𝐛`$. Therefore, $`\lambda _2`$ can be set to unity if all magnetic fields are considered as effective or rescaled magnetic fields. Since the Induction Eq. (4 is linear in the magnetic fields $`𝐛(𝐱,𝐭)`$, such rescaling leaves every conclusion unchanged. Of course, the external force $`𝐠`$ is also scaled by a factor $`\sqrt{\alpha }`$ which does not affect out analysis here as the assignment of canonical dimensions to various fields and parameters is done after absorbing $`\lambda _2`$ in the definition of $`𝐛`$. More specifically under the rescaling $`𝐱l𝐱,tl^zt,u_il^{\chi _u}u_i,b_il^{\chi _b}b_i`$ the bare parameters scale as $`\lambda _{1,3}(l)l^{\chi _u+z1}\lambda _{1,3},\lambda _2(l)l^{2\chi _b\chi _u+z1}\lambda _2`$ and $`\stackrel{~}{B}_ol^{\chi _b\chi _u+z1}\stackrel{~}{B}_o`$. Since there are no fluctuation corrections to the nonlinearities $`\lambda _1,\lambda _2,\lambda _3`$, which are the consequences of the invariances under the transformations TI and TII, they can be kept invariant under rescaling of space and time as mentioned above leading to $`\chi _u=\chi _b=\chi `$ and $`\chi +z=1`$. Thus, under naive rescaling the bare parameter $`\stackrel{~}{B}_o\stackrel{~}{B}_ol^{z1}`$. Furthermore, the invariance under the transformation TII and the resulting Ward identity ensure that different ways of implementing one-loop RG by having different values for $`𝐛`$ and $`\stackrel{~}{𝐁}_𝐨`$ in Eqs. (9) and (10) subject to the same bare mean magnetic field $`𝐁_𝐨=\stackrel{~}{𝐁}_𝐨+𝐛`$. Since any renormalisation scheme must respect a freedom of choice as represented by the transformation TII, respective terms must scale the same way under the rescaling of space and time abpassive . Hence, due to fluctuation corrections $`\stackrel{~}{B}_o(l)l^\chi `$. Furthermore, since both $`\stackrel{~}{𝐁}_𝐨`$ and $`𝐛`$ have the same scaling dimensions given by the exponent $`\chi `$, the physical mean magnetic field in the system $`𝐁_𝐨=\stackrel{~}{𝐁}_𝐨+𝐛(𝐱,𝐭)`$, if it receives fluctuation corrections under mode eliminations, must scale as $`B_o(l)l^\chi `$, or if there are no fluctuation corrections to it due to some special symmetries (as in the one-dimensional Burgers-like model for MHD in Ref.abjkb ) will be irrelevant (in an RG) sense as it will flow to zero as $`l^{z1}`$ (since $`z`$ for fully developed turbulence is less than unity). This clearly suggests the possibility of a renormalisation of the Alfvén wave speed, akin to the renormalisation of the sound speed in compressible fluid turbulence jkbsound . As we will see below, for $`3d`$MHD there are infra-red singular fluctuation corrections to the bare mean magnetic field leading it to renormalise in a way consistent with our predictions from the Ward identities above. It should be noted that our analyses above is independent of the strength of the bare value of the mean magnetic field. Although, in the discussion above we have rescaled space isotropically (i.e., $`x,y`$ and $`z`$ coordinates are rescaled the same way) and analyze the scale dependence of the resulting effective parameters and the fields, it should be noted that such a rescaling does not imply that the effective parameters have isotropic structures.
## V Energy spectra for incompressible fluids
We begin by writing down the $`3d`$MHD equations in the incompressible limit (i.e., $`𝐮=\mathrm{𝟎}`$). We write down the equations of motion in $`𝐤`$-space in terms of the Elsässer variables $`𝐳^\pm =𝐮\pm 𝐛`$. The equations are (we take the mean magnetic field $`𝐁_𝐨`$ to be along the $`\widehat{z}`$ direction)
$`{\displaystyle \frac{z_l^+(𝐤,t)}{t}}iB_ok_zz_i^++iP_{lp}k_s{\displaystyle z_s^{}(𝐪,t)z_p^+(𝐤𝐪,t)}+\eta _+^ok^2z_l^++\eta _{}^ok^2z_l^{}`$ $`=`$ $`\theta _l^+(𝐤,t),`$
$`{\displaystyle \frac{z_l^{}(𝐤,t)}{t}}+iB_ok_zz_i^{}+iP_{lp}k_s{\displaystyle z_s^+(𝐪,t)z_p^{}(𝐤𝐪,t)}+\eta _+^ok^2z_l^{}+\eta _{}^ok^2z_l^+`$ $`=`$ $`\theta _l^{}(𝐤,t).`$ (11)
The stochastic forces $`\theta _l^\pm `$ are linear combinations of $`f_l,g_l`$, the tensor $`P_{lp}`$ is the transverse projection operator which appears due to the divergence-free conditions on $`𝐮`$ and $`𝐛`$, and $`\eta _\pm ^o=\nu \pm \eta `$.
In the absence of any crosscorrelations between the velocity and the magnetic fields, the simplest choice for the noise variances, consistent with the divergence-free conditions on the velocity and the magnetic fields, are
$`\theta _l^\pm (𝐤,t)\theta _m^\pm (𝐤,t)=2P_{lm}D_1k^y\delta (t),`$
$`\theta _l^\pm (𝐤,t)\theta _m^{}(𝐤,t)=2P_{lm}D_2k^y\delta (t).`$ (12)
In Eqs.(12) we choose $`y>0`$. In particular, the choice of $`y=d`$ in $`d`$-space dimensions ensures that the energy flux in the inertial range is a constant (up to a weak logarithmic dependence on wavevector $`k`$) which forms the basis of the K41 theory k41 , and as a result, the parameters $`D_1`$ and $`D_2`$ pick up dimensions of energy dissipation rate per unit mass. Hence, we will use $`y=d`$ in our calculations below. In experimental realisations of MHD, external forces act on the large scales. In analytical approaches, such forces are replaced by stochastic forces with variances given by Eqs. (12) for calculational convenience. It should however be noted that in numerical solutions of the Navier-Stokes equation driven by stochastic forces with variances similar to Eqs. (12) structure functions of the velocity fields are shown to exhibit multiscaling similar to the experimental results rfnse . Our preliminary results from the numerical solutions of the isotropic $`3d`$MHD equations (i.e., no mean magnetic fields) yield multiscaling similar to those obtained from the $`3d`$MHD Eqs. driven by large-scale forces rfmhd . Although no such numerical results exist for the case with a mean magnetic field, it will presumably be true for such a situation also. In a stochastically driven Langevin description, correlation functions are proportional to the noise correlations and hence it is necessary to force the Induction Eq. (6) stochastically, as is common in the literature (see, e.g., abjkb ; abepl and references therein). For such noises in the absence of any mean magnetic field (i.e., for the isotropic case) it has been shown by using a one-loop mode coupling theory that the scaling exponents have values $`\chi =1/3,z=2/3`$ abepl .
The Eqs. of motion (11) are coupled even at the linear level leading to a bare propagator matrix for the Eqs. (11) of the form (a ’$`o`$’ refers to bare quantities)
$$G_{o}^{}{}_{}{}^{1}=\left(\begin{array}{cc}i\omega iB_ok_z+\eta _+^ok^2& \eta _{}^ok^2\\ \eta _{}^ok^2& i\omega +iB_ok_z+\eta _+^ok^2\end{array}\right)$$
We use a one-loop mode coupling theory which is conveniently formulated in terms of the self-energy matrix $`\mathrm{\Sigma }`$ and the correlation functions given by Eqs. (7). The self-energy matrix $`\mathrm{\Sigma }`$ is defined by $`G^1=G_o^1\mathrm{\Sigma }`$ where $`G`$ is the renormalised propagator matrix. In terms of the scaling exponents $`\chi `$ and $`z`$, and the renormalized parameters the self-energy matrix $`\mathrm{\Sigma }`$ is given by
$$\mathrm{\Sigma }=\left(\begin{array}{cc}i\omega iB(k)k_z+\eta _+k^z& \eta _{}k^z\\ \eta _{}k^z& i\omega +iB(k)k_z+\eta _+k^z\end{array}\right)$$
where $`B(k)`$ is the renormalised or the effective Alfvén wave speed. If there are diagrammatic corrections to the imaginary parts of $`\mathrm{\Sigma }(k,\omega )`$ at frequency $`\omega =0`$ which are singular in the infra-red limit, then $`B(k)`$ is different from $`B_o`$, the bare Alfvén wave speed, else they are the same.
There are two one-loop diagrams which contribute to the fluctuation corrections to $`\mathrm{\Sigma }_{11}(k,\omega =0)=iB(k)k_z+\eta _+k^z`$. These are shown in Fig.(1). These have both real and imaginary parts at frequency $`\omega =0`$. Thus there are fluctuation corrections to the bare Alfvén wave speed which are singular in the small wavenumber limit. We assume, in the spirit of mode-coupling methods, $`B(k)=Bk^s`$ (we assume an isotropic scale dependence for $`B(k)`$ for simplicity which suffices our purpose of finding the scale dependence of the effective Alfvén wave speed). Clearly, if $`s+1<z`$, the small wavenumber limit of the problem is dominated by underdamped waves, if $`s+1>z`$ the Alfvén waves are damped out in the small wavenumber limit. In contrast if $`s+1=z`$ then in the small wavenumber limit both the propagating and the dissipative modes are present. Our analyses in Sec.IV clearly suggest that if there are fluctuation corrections to $`B_o`$ then the effective scale dependent Alfvén wave speed $`B(l)k^\chi `$ yielding $`s=\chi `$. Furthermore, since $`\chi +z=1`$ we have $`s+1=z`$ leading to the co-existence of the underdamped Alfvén waves and the dissipative modes in the large-scale, long-time limit. Note that this situation allows us to define a dimensionless parameter $`\beta B/\eta _+`$ where $`B`$ and $`\eta _+`$ are the renormalised Alfvén speed and the viscosity respectively.
In the SCMC approach vertex corrections are neglected. Lack of vertex renormalisations in the zero wavevector limit in $`3d`$MHD allows SCMC to yield exact relations between the scaling exponents $`\chi `$ and $`z`$ as in the noisy Burgers/Kardar-Parisi-Zhang equation freykpz . In this model the nonlinearities and the noise variances do not renormalise, thus leading to $`z=2/3,\chi =1/3`$, satisfying the exponent identity $`\chi +z=1`$. In the present case, due to the singular nature (in the small wavevector limit) of the bare noise correlations (12), they do not pick up any further singular corrections to their scaling; however the amplitudes get modified (this is similar to the results in Ref.abepl ). We denote the renormalized amplitudes by $`D`$ and $`\stackrel{~}{D}`$, respectively. Furthermore, for the $`3d`$MHD Eqs. (5) and (6) there are infra-red singular fluctuation corrections to $`B_o`$ \[see the one-loop diagrams in Fig. (1)\] leading to $`B_o(l)l^\chi `$ consistent with our arguments above. The SCMC approach involves consistency in the scaling and the amplitudes of the mode coupling equations. It should be noted that the exponent values $`z=2/3,\chi =1/3`$ satisfy the mode coupling integral equations regardless of the strength of the mean magnetic field. For amplitude consistency the one-loop integrals are required to be evaluated. Due their complicated structure, we evaluated them by assuming that the strength of the mean magnetic field is small. Therefore, only our amplitude relations and not the scaling exponents are affected by the approximation of small mean magnetic fields.
Demanding amplitude consistency in the mode coupling equations we obtain,
$`\eta _+`$ $`=`$ $`{\displaystyle \frac{D}{\eta _+^2}}\left[1{\displaystyle \frac{2}{d}}+{\displaystyle \frac{1}{d(d+2)}}\right]+{\displaystyle \frac{\stackrel{~}{D}}{d(d+2)\eta _+^2}}{\displaystyle \frac{\beta ^2D}{\eta _+^2d}}\left[1{\displaystyle \frac{2}{d}}+{\displaystyle \frac{1}{d(d+2)}}{\displaystyle \frac{\beta ^2\stackrel{~}{D}}{\eta _+^2d(d+2)(d+4)}}\right],`$
$`\beta `$ $``$ $`{\displaystyle \frac{B}{\eta ^+}},B=B_o{\displaystyle \frac{D}{d(d+2)\eta _+^3}},`$
$`{\displaystyle \frac{1\mathrm{\Gamma }}{1+\mathrm{\Gamma }}}`$ $`=`$ $`{\displaystyle \frac{(1+\mathrm{\Gamma }^2)(1\frac{2}{d}+\frac{2}{d(d+2)})0.5\mathrm{\Gamma }(1\frac{2}{d})+F_1(\beta )}{(1+\mathrm{\Gamma }^2)(1\frac{2}{d}+\frac{2}{d(d+2)})+0.5\mathrm{\Gamma }(1\frac{2}{d})+F_2(\beta )}},`$ (13)
where $`F_1(\beta )2\beta ^2(1+3\mathrm{\Gamma }^2)(1/d+\frac{1}{d(d+2)}\frac{1}{d(d+2)(d+4)}),F_2(\beta )2\beta ^2(3+\mathrm{\Gamma }^2)(1/d+\frac{1}{d(d+2)}\frac{1}{d(d+2)(d+4)})`$ are the $`\beta `$-dependent parts of the amplitude-ratio. While calculating the amplitude consistency relations we worked in the limit of small $`\beta B/\eta _+`$ and the renormalised magnetic Prandtl number ($`\frac{\eta }{\nu }=\frac{\eta _+\eta _{}}{\eta _++\eta _{}}`$ ) close to unity. Away from these limits the amplitudes of the underlying one-loop integrals become much more complicated functions of $`\beta `$ and also of the renormalised magnetic Prandtl number, but the qualitative picture remains unchanged. The main physical picture that emerges from the expressions (13) is that in the absence of a bare mean magnetic field $`B_o`$ the effective Alfvén wave speed is also zero which is a restatement of the fact that if the original theory is isotropic, so will be the renormalised theory. Moreover the renormalised Alfvén speed increases with increasing $`B_o`$. With these results we are now in a position to calculate the energy spectra in the inertial range. We use the form of the effective (scale dependent) viscosity and the Alfvén wave speed to obtain the equal time correlation function of $`z_i^\pm `$ in the long wavelength limit. We define (in $`3d`$)
$`C_{ij}^{++}(k,\omega )`$ $`=`$ $`z_i^+(k,\omega )z_j^+(k,\omega )={\displaystyle \frac{2Dk^3P_{ij}(𝐤)}{(\omega Bk^{1/3}k_z)^2+\eta _+^2k^{4/3}}},`$
$`C_{ij}^{}(k,\omega )`$ $`=`$ $`z_i^{}(k,\omega )z_j^{}(k,\omega )={\displaystyle \frac{2Dk^3P_{ij}(𝐤)}{(\omega +Bk^{1/3}k_z)^2+\eta _+^2k^{4/3}}},`$
$`C_{ij}^+(k,\omega )`$ $`=`$ $`z_i^+(k,\omega )z_j^{}(k,\omega )={\displaystyle \frac{2\stackrel{~}{D}k^{4/3}P_{ij}(𝐤)}{[i(\omega Bk^{1/3}k_z)+\eta _+k^{2/3}][i(\omega +Bk^{1/3}k_z)+\eta _+k^{2/3}]}}.`$ (14)
As discussed before, we have omitted anisotropic corrections of $`O(B_o)^2`$ or of $`O(\beta )^2`$ to the correlation functions as we are trying to find out the relevance (in an RG sense) of Alfvén waves on the scaling of the isotropic correlation functions in the absence of any mean magnetic fields. Therefore, the equal time correlation functions have the following form (in $`d=3`$):
$$C_{ij}^{++}(k,t=0)=C_{ij}^{}(k,t=0)=Dk^{32/3}P_{ij}(𝐤).$$
(15)
Therefore, the equal time autocorrelation functions of $`𝐳^\pm `$ are independent of any mean magnetic field. This holds true regardless of the scaling of the noise variances. The equal time cross correlation function $`C_{ij}^+(k,t=0)`$ requires more careful considerations. On integrating $`C_{ij}^+(k,\omega )`$ over all frequency $`\omega `$ one obtains (in $`3d`$)
$$C_{ij}^+(k,t=0)=\frac{\stackrel{~}{D}k^3P_{ij}}{iBk^{1/3}k_z+\eta _+k^{2/3}},$$
(16)
a form which is valid in the inertial range. It is clear from the expression (16) that both the real and the imaginary parts of $`C_{ij}^+(k,t=0)`$ scale as $`k^{32/3}`$ in the inertial range. Thus the one-dimensional kinetic- and the magnetic-energy spectra (which are simply related to the correlators defined in Eqs. (14)) scale as $`k^{5/3}`$ in the inertial range. The emerging physical picture is as follows: We find from the expression for $`C_{ij}^+(k,\omega )`$ above that this, as a function of frequency $`\omega `$, has maxima at $`\omega =\pm Bk^{1/3}k_z`$ and the width at half-maxima $`k^{2/3}`$. In contrast, the auto correlations of $`𝐳^+,𝐳^{}`$ are maximum at $`\omega =0`$ and their widths scale as $`k^{2/3}`$. Thus in the long wavelength limit the width and the location of the maxima of $`C_{ij}^+(k,\omega )`$ scale in the same way leading to the presence of the underdamped Alfvén waves in the hydrodynamic limit. Therefore, it immediately follows that the kinetic- and the magnetic-energy spectra, being linear combinations of the correlators discussed above times appropriate phase factors, scale as $`k^{5/3}`$ even in the presence of a non-zero mean magnetic field. It should be noted that we have used a small $`\beta `$ approximation to arrive at our results for the self-consistent amplitude relations. For a finite $`\beta `$ one would require to work with a fully anisotropic form of the correlation functions and obtain self-consistent relations for scaling and anisotropic amplitudes, a task, which is analytically challenging, remains to be done in the future. However, based on our calculations above, the exponent identity $`\chi +z=1`$ and the Ward identity suggesting that the renormalised Alfvén wave speed should scale as $`k^{1/3}`$ with wavevector $`k`$ being in the inertial range, we argue that even for a finite $`\beta `$, i.e., a finite mean magnetic field the energy spectra will scale as $`k^{5/3}`$ in the inertial range, a result supported by many observational evidences obser . The self-consistent amplitude equations will then be anisotropic reflecting the presence of underdamped Alfvén waves in the inertial range. The full correlation matrix will be anisotropic in the hydrodynamic limit; its eigenvalues have different amplitudes, but all of them scale the same way. Our confidence on our result that the scaling of the correlation functions along directions parallel and perpendicular to the direction of the mean magnetic field is same is derived from the fact that our exponent values $`z=2/3,\chi =1/3`$ satisfy the one-loop integral equations regardless of the strength of the mean magnetic field and are consistent with the Ward identity discussed above.
## VI Kolmogorov’s constants
According to the Kolmogorov’s hypothesis for fluid turbulence k41 , in the inertial range energy spectrum $`E(k)=K_oϵ^{2/3}k^{5/3}`$, where $`K_o`$, a universal constant, is the Kolmogorov’s constant and $`ϵ`$ is the energy dissipation rate per unit mass. Various calculations, based on different techniques by different groups krai1 ; leslie ; yakhot ; jkbamita show that $`K_o1.5`$ in three dimensions. Having noted that the energy spectra, even in the presence of a mean magnetic field scale as $`k^{5/3}`$ extensions of Kolmogorov’s hypothesis for $`3d`$MHD allows one to define Kolmogorov’s constants for the Elsässer fields: $`E_\pm (k)=K_o^\pm ϵ_\pm ^{2/3}k^{5/3}`$. Since $`𝐳^\pm =𝐮\pm 𝐛`$, we have $`z_i^+(𝐤,\omega )z_j^+(𝐤,\omega )=z_i^{}(𝐤,\omega )z_j^{}(𝐤,\omega )`$ and $`ϵ_+=ϵ_{}=ϵ_{MHD}`$ in absence of any crosscorrelations between the velocity and the magnetic fields, we have $`K_o^+=K_o^{}=K_{MHD}`$. The noise strength $`D`$ and the rate of energy dissipation per unit mass is connected by the Novikov’s theorem nov ; abepl :$`ϵ=2D\frac{S_3}{(2\pi )^3}`$. Noting that the energy spectra $`E_\pm (k)`$ of $`𝐳^\pm `$ in the inertial range is given by
$$E_\pm (k)=Dk^3/\eta _+k^{2/3},$$
(17)
where $`D`$ is the effective or renormalised noise strength, we identify
$$K_{MHD}=1.6\left[1+0.7\left(3\mathrm{\Gamma }^26\mathrm{\Gamma }\beta ^229/105\right)\right]^{2/3}.$$
(18)
The notable feature of the expression (18) is that the constant $`K^{MHD}`$ depends on the dimensionless parameter $`\beta `$ which we introduced before. For $`\mathrm{\Gamma }=0`$, i.e., for no magnetic fields, we find $`K_{MHD}=K_o=1.6`$ for pure fluid turbulence which is well-within the accepted range of values leslie . Before closing this Sec. we would like to point out that the presence of multiscaling raises questions about Kolmogorov’s constant $`K_{MHD}`$ being universal: A small but finite intermittency correction (i.e., multiscaling) over the simple K41 scaling implies the presence of an arbitrary scale which may spoil the universality of $`K_{MHD}`$. We however refrain ourselves from getting into this question and adopt a point of view that regardless of whether or not $`K_{MHD}`$ remains universal due to multiscaling, the numerical value of this constant is likely to get affected by the presence of Alfvén waves in the system which is reflected by the expression (18).
## VII Possibilities of variable multifractality
Experiments and numerical sumulations abprl ; cho find nonlinear multiscaling corrections to the K41 prediction of $`\zeta _p^a=p/3`$ for the structure functions in the inertial range. Until the date, no controlled perturbative calculation for $`\zeta _p^a`$ is available. To account for multiscaling in fluid turbulence, however, Obukhov obu and Kolmogorov kol2 assumed a log-normal distribution for dissipation $`ϵ`$ to arrive at
$$S_p^v(r)=|\mathrm{\Delta }v|^p=C_p\overline{ϵ}^{p/3}r^{p/3}\left(\frac{L}{r}\right)^{(\overline{\delta }/2)p(p3)},$$
(19)
where $`\overline{ϵ}`$ is the mean value of $`ϵ`$ and (a bound for $`ϵ`$ in fluid turbulence has been discussed in Ref.doer ):
$$ϵ(𝐱+𝐫)ϵ(𝐱)(\mathrm{\Delta }v)^6/r^2(L/r)^{9\overline{\delta }}.$$
(20)
For small $`\overline{\delta }`$, $`\delta 9\overline{\delta }`$. A standard calculation on the randomly stirred model yields intermittency exponent $`\delta =0.2`$ jkbamita where $`\delta =9\overline{\delta }`$, whereas the best possible estimate from experiments is 0.23 jkbamita . This model, despite having well-known limitations and difficulties frisch , serves as a qualitative illustration of multiscaling. As in fluid turbulence, in MHD the dissipations $`ϵ_\pm `$ of the Elsässer variables $`𝐳^\pm `$ fluctuate in space and time, and as a result one may define two intermittency exponents $`\delta _\pm `$ for them. In the present problem $`\delta _+=\delta _{}=\delta _{MHD}`$. Below we calculate the exponent $`\delta _{MHD}`$ in a one-loop expression which will give us an estimate of the Alfvén wave speed-dependent deviation of the scaling of the structure functions from their simple-scaling values as predicted by the K41 theory. We closely follow Ref.jkbamita . We work with the self-consistent forms for the self-energies and correlation functions given above along with the consistency relations for the amplitude-ratios $`\mathrm{\Gamma }\mathrm{and}\beta `$. Following Ref.jkbamita , we find the dissipation correlation functions in $`3d`$ to be
$$ϵ(𝐱+𝐫)ϵ(𝐱)12.4ϵ_{MHD}^2\alpha ^2K_{MHD}^2\mathrm{ln}\frac{L}{r},$$
(21)
with $`\alpha `$ being defined by the relation $`\nu _+=\alpha ϵ_{MHD}^{1/3}`$. From the self-consistent amplitude-relations (13) we find
$$\alpha =0.4\left[1+0.7\left(3\mathrm{\Gamma }^26\mathrm{\Gamma }\beta ^229/105\right)\right]^{1/3}\left[10.5\mathrm{\Gamma }4\beta ^2(0.70.03\mathrm{\Gamma })\right]^{1/3}.$$
(22)
Thus, $`\alpha `$ which is a universal coefficient in ordinary fluid turbulence, varies with the parameter $`\beta `$, or with the Alfvén wave speed in MHD. Substituting the values of $`K_o`$ and $`\alpha `$ we find
$$\delta _+=\delta _{}=\delta _{MHD}=0.2\left[1+0.7\left(3\mathrm{\Gamma }^26\mathrm{\Gamma }\beta ^229/105\right)\right]^{4/3}.$$
(23)
Thus we find that a decrease in the value of $`\delta `$ with an increase in $`\beta `$, i.e., with increasing mean magnetic field. Despite the limited applicabilities of log-normal models in characterising multiscaling, we can conclude, from our expression (23) for the intermittency exponent in MHD in the presence of a mean magnetic field, that the intermittency corrections to the simple K41 scaling is likely to get affected (reduced in our calculations) in the presence of Alfvén waves. Further calculations and/or numerical simulations are needed to find the exact extent of the dependence of multiscaling on the mean magnetic field and the possibilities of anisotropic multiscaling for the structure functions parallel and perpendicular to the mean magnetic field as demonstrated recently in Ref.biskampnew . It should be noted that our conclusions on the multiscaling properties of the structure functions depend on a log-normal model for $`3d`$MHD. Such a description, unfortunately, is unable to distinguish between the structure functions parallel and perpendicular to the direction of the mean magnetic field in the system. Moreover, the intermittency exponents above \[expressions (23)\] are evaluated in the lowest order in mean magnetic field. Therefore, even though from our results we are not able to make firm comments on the possibility of parallel and perpendicular structure functions exhibiting different multiscaling, the real importance of our results lie in their elucidation of the multiscaling properties depending on the mean magnetic field.
## VIII Conclusion
In this paper we have considered the effects of the Alfvén waves on the statistical properties of the correlation functions of the velocity and the magnetic fields or the Elsässer fields. We considered the case when the velocity fields are incompressible. In a one-loop approximation we find that the effective or the renormalised Alfvén wave speed scales as $`k^{1/3}`$ where the wavevector $`k`$ is in the inertial range. This immediately yields that the energy spectra, even in the presence of a mean magnetic field, scale as $`k^{5/3}`$ in the inertial range. We identify a dimensionless parameter $`\beta `$ which is the ratio of the effective Alfvén wave speed and the renormalised viscosity. We obtain self-consistent relations between the amplitude ratio $`\mathrm{\Gamma }`$ of the correlation functions and $`\beta `$. These relations allow us to calculate the dimensionless Kolmogorov’s constant and we show that it depends explicitly on $`\beta `$ or on the mean magnetic field. Finally, we calculate the intermittency exponent which in a log-normal model gives a qualitative account of the multiscaling in terms of the deviation from the K41 scaling for the structure functions. We would like to emphasize that although the one-dimensional Burgers-like model of MHD of Ref.abjkb and the $`3d`$MHD Eqs. yield energy spectra independent of the Alfvén waves, the long wavelength physical pictures are different. In the former case, due to the nonrenormalization of the bare Alfvén wave speed, Alfvén waves are overdamped in the hydrodynamic limit; the dominant process in that limit is viscous dissipation. In contrast, in $`3d`$MHD, the Alfvén wave speed picks up a singular correction in the hydrodynamic limit leading to the K41 scaling of the energy spectra and the presence of underdamped Alfvén waves in the hydrodynamic limit. Despite similar mathematical structures these crucial differences arise principally because the one-dimensional model decouples completely when written in terms of the Elsässer variables, allowing to comove with the waves of each of them separately. As a result correlators of each of them are independent of the Alfvén waves and the bare Alfvén wave speed remains unrenormalised leading to overdamped Alfvén waves in the hydrodynamic limit. However, the $`3d`$MHD equations do not decouple when written in terms of the Elsässer variables and hence oppositely propagating waves cannot be made to vanish by comoving. Despite the limitations of the one-loop methods yakhot and the small $`\beta `$ approximation to facilitate easier analytical manipulations for the amplitude relations, we obtain results which are significantly new and open the intriguing possibilities of MHD multiscaling universality classes being parametrised by the mean magnetic field. Some of the quantitative details will change if one retains terms which are higher order in $`\beta `$; however, we believe that the qualitative picture will essentially remain the same. As MHD turbulence forms a natural example of a driven nonequilibrium system with Alfvén waves, our results are the first of its kind for a natural system. In our log-normal model approach, we did not distinguish between the longitudinal and the transverse structure functions our results cannot be compared directly with those of Ref.biskampnew , where the longitudinal and the transverse structure functions are found to scale differently with multiscaling exponents which depend on the magnitude of the mean magnetic field in the presence of a finite magnetic helicity. In contrast our results apply to the multiscaling of the usual structure functions in the absence of any magnetic helicity which are combinations of the transverse and the longitudinal structure functions as considered in Ref.biskampnew which would also then exhibit a mean magnetic field dependent scaling. In particular we find that the usual structure functions multiscale less in presence of an increasingly strong mean magnetic field which is in qualitative agreement with those of Ref.biskampnew . Although here we have restricted ourselves to the study of Alfvén waves as relevant perturbations on the amplitude and the scaling of the isotropic correlation functions in the limit of small $`\beta `$, our results indicate the possibility of the multiscaling exponents varying with the amplitude of the mean magnetic field. Our analyses, when extended for finite values of $`\beta `$ and for full anisotropic structure of the correlation functions, are likely to provide further understanding of and resolve some of the discrepancies between the various phenomenological scenarios available for situations with large mean magnetic fields galtier ; ng ; montgomery . We leave this task for the future.
From a broader point of view our results demonstrate the critical role of wave-like excitations in determining the statistical properties of $`3d`$MHD. The presence of propagating waves is not confined to MHD only; they are a generic feature in many other naturally occurring soft-matter systems where such waves can be present in a viscous environment, e.g., active polar gels in cytoskeletal dynamics frank and in the dynamics of self-propelled particles sriram etc. We believe our results will lead to similar theoretical and experimental studies in relevant nonequilibrium soft-matter systems.
## IX Acknowledgement
This work was started when one of the authors (AB) was an Alexander von Humboldt Fellow at the Hahn Meitner Institut, Berlin. AB wishes to thank the AvH Stiftung, Germany for partial financial support. |
warning/0506/hep-lat0506011.html | ar5iv | text | # DESY/05-090SFB/CPP-05-20June 2005 HMC algorithm with multiple time scale integration and mass preconditioning
## 1 Introduction
Simulations of full QCD with light dynamical flavors of quarks and at small lattice spacing are presently one of the greatest challenges in lattice QCD. Simulations in this regime have to face the phenomenon of critical slowing down: in addition to the “natural” increase of costs due to increasing volume and increasing iteration numbers needed in the solvers, the autocorrelation times are expected to grow significantly when the masses are decreased.
One widely used algorithm to perform those simulations is the Hybrid Monte Carlo (HMC) algorithm , an exact algorithm which combines molecular dynamics evolution of the gauge fields with a Metropolis accept/reject step to correct for discretization errors in the numerical integration of the corresponding equations of motion. However, in its original form the HMC algorithm is even on nowadays computers not able to tackle simulations with light quarks on fine lattices, for the reasons stated above.
Therefore a lot of effort has been invested to accelerate the HMC algorithm during the last years and the corresponding list of improvements that were found is long. It reaches, for example, from even/odd preconditioning over multiple time scale integration and chronological inversion methods to mass preconditioning (Hasenbusch acceleration) , to mention only those that are immediately relevant for the present paper. It is worth noting that many of the known improvement tricks can be combined. In addition, alternative multiboson methods have been suggested, which, however, appear not to be superior to the HMC algorithm.
Recently in ref. a HMC variant as a combination of multiple time scale integration with domain decomposition as preconditioner on top of even/odd preconditioning was presented and speedup factors of about $`10`$ were reported when compared to state of the art simulations using a HMC algorithm . Even more important, excellent scaling properties with the quark mass were found. Thus this algorithm seems to be most promising when one wants to simulate small quark masses on fine lattices.
In this paper we are going to present yet another variant of the HMC algorithm similar to the one of refs. comprising multiple time scale integration with mass preconditioning on top of even/odd preconditioning. We test this algorithm for standard Wilson fermions at $`\beta =5.6`$ and at pion masses ranging from $`m_\pi =370`$MeV to $`m_\pi =650`$MeV. We show that in this situation the algorithm has similar scaling properties and performance as the method presented in ref. . From the performance data obtained with our HMC variant we update the “Berlin Wall” figures of refs. and find that the picture is significantly improved.
After shortly recalling the HMC algorithm as well as the ideas of multiple time scale integration and mass preconditioning, we will present numerical results comparing our HMC version with the one of refs. and draw our conclusions.
## 2 HMC algorithm
The variant of the HMC algorithm we will present here is applicable to a wide class of lattice Dirac operators, including twisted mass fermions, various improved versions of Wilson fermions, staggered fermions, and even the overlap operator. Nevertheless, in order to discuss a concrete example, we restrict ourselves in this paper to the Wilson-Dirac operator with Wilson parameter $`r`$ set to one
$$D_\mathrm{W}[U,m_0]=\frac{1}{2}\underset{\mu }{}\left\{(_\mu +_\mu ^{})\gamma _\mu a_\mu ^{}_\mu \right\}+m_0,$$
(2-1)
where $`m_0`$ is the bare mass parameter, and $`_\mu `$ and $`_\mu ^{}`$ the gauge covariant lattice forward and backward difference operators, respectively. The Wilson lattice action can be written as
$$S[U,m_0]=S_\mathrm{G}[U]+a^4\underset{x}{}\overline{\psi }(x)\left(D_\mathrm{W}[U,m_0]\right)\psi (x),$$
(2-2)
where $`S_\mathrm{G}[U]`$ is the usual Wilson plaquette gauge action. In the implementation we use the so-called hopping parameter representation of eq. (2-2), where the hopping parameter is defined to be
$$\kappa =(2m_0+8)^1,$$
and the fermion fields are rescaled according to
$$\psi \frac{\sqrt{2\kappa }}{a^{3/2}}\psi ,\overline{\psi }\frac{\sqrt{2\kappa }}{a^{3/2}}\overline{\psi }.$$
After integrating out the fermion fields $`\psi ,\overline{\psi }`$ the partition function of lattice QCD for $`n_f=2`$ mass degenerate flavors of quarks reads
$$𝒵=𝒟Udet(D_\mathrm{W}[U,m_0])^2e^{S_\mathrm{G}[U]}.$$
(2-3)
Since $`D_\mathrm{W}`$ fulfills the property $`\gamma _5D_\mathrm{W}\gamma _5=D_\mathrm{W}^{}`$ we define for later purpose the hermitian Wilson-Dirac operator
$$Q=\gamma _5D_\mathrm{W}.$$
(2-4)
In order to prepare the set up for the Hybrid Monte Carlo (HMC) algorithm , the determinant $`det(D_\mathrm{W}[U,m_0])`$ is usually re-expressed by the use of so-called pseudo fermion fields $`\varphi `$
$$det(D_\mathrm{W})^2=det(Q)^2𝒟\varphi ^{}𝒟\varphi \mathrm{exp}\left(S_{\mathrm{PF}}[U,\varphi ^{},\varphi ]\right),$$
(2-5)
where $`S_{\mathrm{PF}}[U,\varphi ^{},\varphi ]=|Q^1\varphi |^2`$ is the pseudo fermion action. The pseudo fermion fields $`\varphi `$ are formally identical to the fermion fields $`\psi `$, but follow the statistic of bosonic fields. The $`\varphi `$ version of the HMC algorithm is based on the Hamiltonian
$$H(P,U,\varphi ,\varphi ^{})=\frac{1}{2}\underset{x,\mu }{}P_{x,\mu }^2+S_\mathrm{G}[U]+S_{\mathrm{PF}}[U,\varphi ,\varphi ^{}],$$
(2-6)
where we introduced traceless Hermitian momenta $`P_{x,\mu }`$ as conjugate fields to the gauge fields $`U_{x,\mu }`$. The HMC algorithm is then composed out of molecular dynamics evolution of the gauge fields and momenta and a Metropolis accept/reject step, which is needed to correct for the discretization errors of the numerical integration of the corresponding equations of motion.
It is possible to prove that the HMC algorithm satisfies the *detailed balance* condition and hence the configurations generated with this algorithm correctly represent the intended ensemble.
### 2.1 Molecular dynamics evolution
In the molecular dynamics part of the HMC algorithm the gauge fields $`U`$ and the momenta $`P`$ need to be evolved in a fictitious computer time $`t`$. With respect to $`t`$, Hamilton’s equations of motion read
$$\frac{dU}{dt}=\frac{dH}{dP}=P,\frac{dP}{dt}=\frac{dH}{dU}=\frac{dS}{dU},$$
(2-7)
where we set $`S=S_\mathrm{G}+S_{\mathrm{PF}}`$ and $`d/dU`$, $`d/dP`$ formally denote the derivative with respect to group elements. Since analytical integration of the former equations of motion is normally not possible, these equations must in general be integrated with a discretized integration scheme that is area preserving and reversible, such as the leap frog algorithm. The discrete update with integration step size $`\mathrm{\Delta }\tau `$ of the gauge field and the momenta can be defined as
$$\begin{array}{cc}\hfill T_\mathrm{U}(\mathrm{\Delta }\tau )& :UU^{}=\mathrm{exp}\left(i\mathrm{\Delta }\tau P\right)U,\hfill \\ \hfill T_\mathrm{S}(\mathrm{\Delta }\tau )& :PP^{}=Pi\mathrm{\Delta }\tau \delta S,\hfill \end{array}$$
(2-8)
where $`\delta S`$ is an element of the Lie algebra of $`\mathrm{SU}(3)`$ and denotes the variation of $`S`$ with respect to the gauge fields. The computation of $`\delta S`$ is the most expensive part in the HMC algorithm since the inversion of the Wilson-Dirac operator is needed. With (2-8) one basic time evolution step of the so called leap frog algorithm reads
$$T=T_\mathrm{S}(\mathrm{\Delta }\tau /2)T_\mathrm{U}(\mathrm{\Delta }\tau )T_\mathrm{S}(\mathrm{\Delta }\tau /2),$$
(2-9)
and a whole trajectory of length $`\tau `$ is achieved by $`N_{\mathrm{MD}}=\tau /\mathrm{\Delta }\tau `$ successive applications of the transformation $`T`$.
### 2.2 Integration with multiple time scales
In order to generalize the leap frog integration scheme we assume in the following that we can bring the Hamiltonian to the form
$$H=\frac{1}{2}\underset{x,\mu }{}P_{x,\mu }^2+\underset{i=0}{\overset{k}{}}S_i[U],$$
(2-10)
with $`k1`$. For instance with $`k=1`$ $`S_0`$ might be identified with the gauge action and $`S_1`$ with the pseudo fermion action of eq. (2-6).
Given a form of the Hamiltonian (2-10) one can think of the following situations in which it might be favorable to use a generalized leap frog integration scheme where the different parts $`S_i`$ are integrated with different step sizes, as proposed in ref. .
Clearly, in order to keep the discretization errors in a leap frog like algorithm small, the time steps have to be small if the driving forces are large. Hence multiple time scale integration is a valuable tool if the forces originating from the single parts in the Hamiltonian (2-10) differ significantly in their absolute values. Then the different parts in the Hamiltonian might be integrated on time scales inverse proportionally deduced from the corresponding forces.
Another situation in which multiple time scales might be useful exists when one part in the Hamiltonian (2-10) is significantly cheaper to update than the others. In this case the cheap part might be integrated without too much performance loss on a smaller time scale, which reduces the discretization errors coming from that part.
The leap frog integration scheme can be generalized to multiple time scales as has been proposed in ref. without loss of reversibility and the area preserving property. The scheme with only one time scale can be recursively extended by starting with the definition
$$T_0=T_{\mathrm{S}_0}(\mathrm{\Delta }\tau _0/2)T_\mathrm{U}(\mathrm{\Delta }\tau _0)T_{\mathrm{S}_0}(\mathrm{\Delta }\tau _0/2),$$
(2-11)
with $`T_\mathrm{U}`$ defined as in eq. (2-8) and where $`T_{\mathrm{S}_i}(\mathrm{\Delta }\tau )`$ is given by
$$T_{\mathrm{S}_i}(\mathrm{\Delta }\tau ):PPi\mathrm{\Delta }\tau \delta S_i[U].$$
(2-12)
As $`\mathrm{\Delta }\tau _0`$ will be the smallest time scale, we can recursively define the basic update steps $`T_i`$, with time scales $`\mathrm{\Delta }\tau _i`$ as
$$T_i=T_{\mathrm{S}_i}(\mathrm{\Delta }\tau _i/2)[T_{i1}]^{N_{i1}}T_{\mathrm{S}_i}(\mathrm{\Delta }\tau _i/2),$$
(2-13)
with integers $`N_i`$ and $`0<ik`$. One full trajectory $`\tau `$ is then composed by $`[T_k]^{N_k}`$. The different time scales $`\mathrm{\Delta }\tau _i`$ in eq. (2-13) must be chosen such that the total number of steps on the $`i`$-th time scale $`N_{\mathrm{MD}_i}`$ times $`\mathrm{\Delta }\tau _i`$ is equal to the trajectory length $`\tau `$ for all $`0ik`$: $`N_{\mathrm{MD}_i}\mathrm{\Delta }\tau _i=\tau `$. This is obviously achieved by setting
$$\mathrm{\Delta }\tau _i=\frac{\tau }{N_kN_{k1}\mathrm{}N_i}=\frac{\tau }{N_{\mathrm{MD}_i}},0ik,$$
(2-14)
where $`N_{\mathrm{MD}_i}=N_kN_{k1}\mathrm{}N_i`$.
In ref. also a partially improved integration scheme with multiple time scales was introduced, which reduces the size of the discretization errors. Again, we assume a Hamiltonian of the form (2-10) with now $`k=1`$. By defining similar to $`T_0`$
$$T_{\mathrm{SW}_0}=T_{\mathrm{S}_0}(\mathrm{\Delta }\tau _0/6)T_\mathrm{U}(\mathrm{\Delta }\tau _0/2)T_{\mathrm{S}_0}(2\mathrm{\Delta }\tau _0/3)T_\mathrm{U}(\mathrm{\Delta }\tau _0/2)T_{\mathrm{S}_0}(\mathrm{\Delta }\tau _0/6),$$
(2-15)
the basic update step of the improved scheme – usually referred to as the Sexton-Weingarten (SW) integration scheme – reads
$$\begin{array}{cc}\hfill T_{\mathrm{SW}_1}=& T_{\mathrm{S}_1}(\mathrm{\Delta }\tau _1/6)\hfill \\ & [T_{\mathrm{SW}_0}]^{N_0}T_{\mathrm{S}_1}(2\mathrm{\Delta }\tau _1/3)\hfill \\ & [T_{\mathrm{SW}_0}]^{N_0}T_{\mathrm{S}_1}(\mathrm{\Delta }\tau _1/6),\hfill \end{array}$$
(2-16)
where $`\mathrm{\Delta }\tau _0=\mathrm{\Delta }\tau _1/(2N_0)`$. This integration scheme not only reduces the size of the discretization errors, but also sets for $`S_0`$ a different time scale than for $`S_1`$. Hence, it is one special example for an integration scheme with multiple time scales and can easily be extended to more than two time scales by recursively defining ($`0<ik`$):
$$\begin{array}{cc}\hfill T_{\mathrm{SW}_i}=& T_{S_i}(\mathrm{\Delta }\tau _i/6)\hfill \\ & [T_{\mathrm{SW}_{i1}}]^{N_{i1}}T_{S_i}(2\mathrm{\Delta }\tau _i/3)\hfill \\ & [T_{\mathrm{SW}_{i1}}]^{N_{i1}}T_{S_i}(\mathrm{\Delta }\tau _i/6).\hfill \end{array}$$
(2-17)
The different time scales for the SW integration scheme are defined by
$$\mathrm{\Delta }\tau _i=\frac{\tau }{(2N_k)(2N_{k1})\mathrm{}(2N_i)}=\frac{\tau }{N_{\mathrm{MD}_i}},ik.$$
(2-18)
Note that the SW partially improved integration scheme makes use of the fact that the computation of the variation of the gauge action is cheap compared to the variation of the pseudo fermion action and in addition the time scales are chosen in order to cancel certain terms in the discretization error exactly.
## 3 Mass Preconditioning
The arguments presented in this section are made for simplicity only for the not even/odd preconditioned Wilson-Dirac operator. The generalization to the even/odd preconditioned case is simple and can be found in ref. and the appendix of ref. .
Mass preconditioning – also known as Hasenbusch acceleration – relies on the observation that one can rewrite the fermion determinant as follows
$$\begin{array}{cc}\hfill det(Q^2)& =det(W^+W^{})\frac{det(Q^2)}{det(W^+W^{})}\hfill \\ & =𝒟\varphi _1^{}𝒟\varphi _1𝒟\varphi _2^{}𝒟\varphi _2e^{\varphi _1^{}\frac{1}{W^+W^{}}\varphi _1\varphi _2^{}W^+\frac{1}{Q^2}W^{}\varphi _2}\hfill \\ & =𝒟\varphi _1^{}𝒟\varphi _1𝒟\varphi _2^{}𝒟\varphi _2e^{S_{\text{PF}_1}S_{\text{PF}_2}}.\hfill \end{array}$$
(3-1)
The preconditioning operators $`W^\pm `$ can in principle be freely chosen, but in order to let the preconditioning work $`W^+W^{}`$ should be a reasonable approximation of $`Q^2`$, which is, however, cheaper to simulate. Moreover, to allow for Monte Carlo simulations, $`det(W^+W^{})`$ must be positive. The generalized Hamiltonian (2-6) corresponding to eq. (3-1) reads
$$H=\frac{1}{2}\underset{x,\mu }{}P_{x,\mu }^2+S_\text{G}[U]+S_{\text{PF}_1}[U,\varphi _1,\varphi _1^{}]+S_{\text{PF}_2}[U,\varphi _2,\varphi _2^{}],$$
(3-2)
and it can of course be extended to more than one additional field.
Note that a similar approach was presented in ref. , in which the introduction of $`n`$ pseudo fermion fields was coupled with the $`n`$-th root of the fermionic kernel.
One particular choice for $`W^\pm `$ is given by
$$W^\pm =Q\pm i\mu ,$$
(3-3)
with mass parameter $`\mu `$ refered to as a twisted mass parameter. In fact the operator
$$\left(\begin{array}{cc}W^+& \\ & W^{}\end{array}\right)=\left(\begin{array}{cc}Q& \\ & Q\end{array}\right)+i\mu \tau _3$$
(3-4)
is the two flavor twisted mass operator with $`\tau _3`$ the third Pauli matrix acting in flavor space. One important property of this choice is that $`W^+W^{}=Q^2+\mu ^2`$. Note that $`(W^+)^{}=W^{}`$ and we remark that in general also $`Q`$ can be a twisted mass operator.
In ref. it was argued that the optimal choice for $`\mu `$ is given by $`\mu ^2=\sqrt{\lambda _{\mathrm{max}}\lambda _{\mathrm{min}}}`$. Here $`\lambda _{\mathrm{max}}`$ ($`\lambda _{\mathrm{min}}`$) is the maximal (minimal) eigenvalue of $`Q^2`$. The reason for the above quoted choice is as follows: the condition number of $`Q^2+\mu ^2`$ is approximately $`\lambda _{\mathrm{max}}/\mu ^2`$ and the one of $`Q^2/(Q^2+\mu ^2)`$ approximately $`\mu ^2/\lambda _{\mathrm{min}}`$. With $`\mu ^2=\sqrt{\lambda _{\mathrm{max}}\lambda _{\mathrm{min}}}`$ these two condition numbers are equal to $`\sqrt{\lambda _{\mathrm{max}}/\lambda _{\mathrm{min}}}`$, both of them being much smaller than the condition number of $`Q^2`$ which is $`\lambda _{\mathrm{max}}/\lambda _{\mathrm{min}}`$.
Since the force contribution in the molecular dynamics evolution is supposed to be proportional to some power of the condition number, the force contribution from the pseudo fermion part in the action is reduced and therefore the step size $`\mathrm{\Delta }\tau `$ can be increased, in practice by about a factor of $`2`$ . Therefore $`Q^2`$ must be inverted only about half as often as before and if the inversion of $`W^+W^{}`$, which is needed to compute $`\delta S_{\mathrm{PF}_1}`$, is cheap compared to the one of $`Q^2`$ the simulation speeds up by about a factor of two .
One might wonder why the reduction of the condition number from $`K`$ to $`\sqrt{K}`$ gives rise to only a speedup factor of about $`2`$. One reason for this is that one cannot make use of the reduced condition number of $`Q^2/(Q^2+\mu ^2)`$ in the inversion of this operator, because in the actual simulation still the badly conditioned operator $`Q^2`$ must be inverted to compute the variation of $`S_{\text{PF}_2}=\varphi _2^{}\frac{W^+W^{}}{Q^2}\varphi _2`$.
### 3.1 Mass preconditioning and multiple time scale integration
In the last subsection we have seen that mass preconditioning is indeed an effective tool to change the condition numbers of the single operators appearing in the factorization (3-1) compared to the original operator. But, this reduction of the condition numbers only influences the forces – which are proportional to some power of the condition numbers of the corresponding operators – and *not* the number of iterations to invert the physical operator $`Q^2`$.
Therefore it might be advantageous to change the point of view: instead of tuning the condition numbers in a way à la refs. we will exploit the possibility of arranging the forces by the help of mass preconditioning with the aim to arrange for a situation in which a multiple time scale integration scheme is favorable, as explained at the beginning of subsection 2.2.
The procedure can be summarized as follows: use mass preconditioning to split the Hamiltonian in different parts. The forces of the single parts should be adjusted by tuning the preconditioning mass parameter $`\mu `$ such that the more expensive the computation of $`\delta S_{\mathrm{PF}_i}`$ is, the less it contributes to the total force. This is possible because the variation of $`(Q^2+\mu ^2)/Q^2`$ is, for $`|\mu |<1`$, reduced by a factor $`\mu ^2`$ compared to the variation of $`1/Q^2`$. In addition, $`W^+W^{}=Q^2+\mu ^2`$ is significantly cheaper to invert than $`Q^2`$. Then integrate the different parts on time scales chosen according to the magnitude of their force contribution.
The idea presented in this paper is very similar to the idea of separating infrared and ultraviolet modes as proposed in ref. . This idea was applied to mass preconditioning by using only two time scales in refs. in the context of clover improved Wilson fermions. However, a comparison of our results presented in the next section to the ones of refs. is not possible, because volume, lattice spacing and masses are different.
## 4 Numerical results
### 4.1 Simulation points
In order to test the HMC variant introduced in the last sections, we decided to compare it with the algorithm proposed and tested in ref. . To this end we performed simulations with the same parameters as have been used in ref. : Wilson-Dirac operator with plaquette gauge action at $`\beta =5.6`$ on $`24^3\times 32`$ lattices. We have three simulation points $`A`$, $`B`$ and $`C`$ with values of the hopping parameter $`\kappa =0.1575`$, $`\kappa =0.1580`$ and $`\kappa =0.15825`$, respectively. The trajectory length was set to $`\tau =0.5`$. The details of the physical parameters corresponding to the different simulation points can be found in table 1. Additionally, this choice of simulation points allows at the two parameter sets $`A`$ and $`B`$ a comparison to results published in ref. , where a HMC algorithm with a plain leap frog integration scheme was used.
### 4.2 Details of the implementation
We have implemented a HMC algorithm for two flavors of mass degenerate quarks with even/odd preconditioning and mass preconditioning with up to three pseudo fermion fields. The boundary conditions are periodic in all directions apart from anti-periodic ones for the fermion fields in time direction. For details of the implementation see the appendix of ref. . For the gauge action the usual Wilson plaquette gauge action is used. The implementation is written in C and uses double precision throughout.
For the mass preconditioning we use
$$W_j^\pm =\gamma _5(D_\mathrm{W}[U,m_0]\pm i\mu _j\gamma _5),$$
(4-1)
with $`j=1,2`$ for the factorization in eq. (3-1), where the $`\mu _j`$ are the additional (unphysical) twisted mass parameters. Therefore, the pseudo fermion actions $`S_{\text{PF}_j}`$ are given by
$$S_{\text{PF}_j}[U]=\{\begin{array}{cc}.\varphi _1^{}\left(\frac{1}{W_1^+W_1^{}}\right)\varphi _1\hfill & j=1.,\hfill \\ .\varphi _j^{}\left(\frac{W_{j1}^+W_{j1}^{}}{Q^2}\right)\varphi _j\hfill & j=N_{\mathrm{PF}}.,\hfill \\ .\varphi _j^{}\left(\frac{W_{j1}^+W_{j1}^{}}{W_j^+W_j^{}}\right)\varphi _j.\hfill & \text{otherwise},\hfill \end{array}$$
(4-2)
where we always chose $`0<\mu _1<\mu _2`$ and $`N_{\mathrm{PF}}`$ denotes the actually used number of pseudo fermion fields.
We have implemented the leap frog (LF) and the Sexton-Weingarten (SW) integration schemes with multiple time scales each as described by eq. (2-13) and eq. (2-17), respectively, where $`k`$ in both equations has to be identified with $`N_{\mathrm{PF}}`$.
The time scales are defined as in eq. (2-14) for the LF integration scheme and as in eq. (2-18) for the SW scheme, with $`N_0`$ corresponding to the gauge action and $`N_j`$ to $`S_{\mathrm{PF}_j}`$ for $`N_{\mathrm{PF}}j>0`$. Note that for the LF integration scheme for one trajectory there are $`N_{N_{\mathrm{PF}}}\mathrm{}N_j+1`$ inversions of the corresponding operator needed, while for the SW integration scheme there are $`2N_{N_{\mathrm{PF}}}\mathrm{}2N_j+1`$ inversions needed.
For the inversions we used the CG and the BiCGstab iterative solvers. We have tested the performance of several iterative solvers for the even/odd preconditioned twisted mass operator with the result, that the CG iterative solver is the best choice in presence of a twisted mass. Thus we used for all inversions of mass preconditioning operators exclusively the CG iterative solver.
For the pure Wilson-Dirac operator $`D_\mathrm{W}`$ the BiCGstab iterative solver is known to perform best . In case of dynamical simulations, however, usually the squared hermitian operator needs to be inverted and in this case the CG is comparable to the BiCGstab. Only in the acceptance step, where $`\gamma _5D_\mathrm{W}`$ (or rather the even/odd preconditioned version of it) needs to be inverted to a high precision, the usage of the CG would be wasteful. For this paper we used the BiCGstab iterative solver for all inversions of either the pure Wilson-Dirac operator itself or $`(\gamma _5D_\mathrm{W})^2`$.
The accuracy in the inversions was set during the computation of $`\delta S_{\mathrm{PF}_j}`$ to $`ϵ_j`$, which means that the inversions were stopped when the approximate solution $`\psi _j`$ of $`A_j\psi _j=\varphi _j`$ fulfills
$$\frac{\varphi _jA_j\psi _j}{\varphi _j}ϵ_j,$$
where $`A_j`$ denotes the operator corresponding to $`S_{\mathrm{PF}_j}`$. During the inversions needed for the acceptance step the accuracy was set to $`\stackrel{~}{ϵ}=10^{10}`$ for all pseudo fermion actions. The inversions in the acceptance step must be rather precise in order not to introduce systematic errors in the simulation, while for the force computation the precision can be relaxed as long as the reversibility violations are not too large. The values of $`ϵ_j`$ and $`\stackrel{~}{ϵ}`$ have been set such that the reversibility violations, which should be under control , are on the same level as reported in ref. , which means that the differences in the Hamiltonian are of the order of $`10^5`$. The values for $`ϵ_j`$ can be found in table 2.
The errors and autocorrelation times were computed with the so called $`\mathrm{\Gamma }`$-method as explained in ref. (see also ref. ), i.e.
$$\tau _{\mathrm{int}}=\frac{1}{2}+\underset{t=1}{\overset{\mathrm{}}{}}\frac{\mathrm{\Gamma }(t)}{\mathrm{\Gamma }(0)},$$
(4-3)
with the autocorrelation function $`\mathrm{\Gamma }(t)`$.
### 4.3 Force contributions
The force contributions to the total force from the separate parts in the action we label by $`F_G`$ for the gauge action and by $`F_j`$ for the pseudo fermion action $`S_{\mathrm{PF}_j}`$. Since the variation of the action with respect to the gauge fields is an element of the Lie algebra of $`\mathrm{SU}(3)`$, we used $`X^2=2\mathrm{Tr}X^2`$ as the definition of the norm of such an element.
In order to better understand the influence of mass preconditioning on the HMC algorithm we computed the average and the maximal norm of the forces $`F_G,F_1,F_2`$ and $`F_3`$ on a given gauge field after all corresponding gauge field updates:
$$\begin{array}{cc}\hfill F_{\mathrm{aver}}& =\frac{1}{4L^3T}\underset{x,\mu }{}F(x,\mu ),\hfill \\ \hfill F_{\mathrm{max}}& =\underset{x,\mu }{\mathrm{max}}\{F(x,\mu )\},\hfill \end{array}$$
(4-4)
and averaged them over all measurements, which we indicate with $`.`$. Examples of force distributions for different runs can be found in figure 1. These investigations lead to the following observations generic to our simulation points:
* With the choice of parameters as given in table 2 the single force contributions are strictly hierarchically ordered with
$`F_G_{\mathrm{aver},\mathrm{max}}>F_1_{\mathrm{aver},\mathrm{max}}>F_2_{\mathrm{aver},\mathrm{max}}>F_3_{\mathrm{aver},\mathrm{max}}`$.
* The maximal force is up to one order of magnitude larger than the average force. This can only be explained by large local fluctuations in this quantity. These fluctuations become larger the smaller the mass is.
Moreover, the force ordering and sizes look very similar to the one reported in ref. .
In a next step we performed some test trajectories without mass preconditioning in order to compare the fermionic forces with and without mass preconditioning. For the value of $`\kappa =0.15825`$ (run $`C`$) the result can be found in figure 2. The bars labeled with $`F`$ correspond to the fermion force without mass preconditioning. The labels $`F_1`$ and $`F_2`$ refer to the two fermionic forces for the run $`C`$ with mass preconditioning. The following ratios are of interest:
$$\begin{array}{cc}& \frac{F_{\mathrm{aver}}}{F_1_{\mathrm{aver}}}1,\frac{F_{\mathrm{aver}}}{F_2_{\mathrm{aver}}}42,\hfill \\ & \frac{F_{\mathrm{max}}}{F_1_{\mathrm{max}}}1.3,\frac{F_{\mathrm{max}}}{F_2_{\mathrm{max}}}29.\hfill \end{array}$$
These ratios show that the average and maximal norm of $`F_2`$ is strongly reduced compared to the average and maximal norm of $`F`$. We observe that the maximal norm is slightly less reduced than the average norm and, by varying $`\mu _1`$, we could confirm that the norm (average and maximal) of $`F_2`$ is roughly proportional to $`\mu _1^2`$.
As a further observation, one sees from figure 2 or from the ratios quoted above that the norm of $`F_1`$ is almost identical to the norm of $`F`$, which is the case for both the average and the maximal values.
From these investigations we think one can conclude the following: in the first place it is possible to tune the value of $`\mu _1`$ (and possibly $`\mu _2`$) such that the most expensive force contribution of $`F_2`$ (or $`F_3`$) to the total force becomes small. Secondly, since in the example above the force contributions for $`F`$ and $`F_1`$ are almost identical – even though the masses are very different – we conclude that the norm of the forces does not explain the whole dynamics of the HMC algorithm. For this point see also the discussion in the forthcoming subsection.
### 4.4 Tuning the algorithm
As mentioned already in section 3.1 the tuning of the different mass parameters and time scales could become a delicate task. Therefore we decided to tune the parameters $`\mu _1`$ and possibly $`\mu _2`$ such that the molecular dynamics steps number $`N_{N_{\mathrm{PF}}}`$ for the LF or $`2N_{N_{\mathrm{PF}}}`$ for the SW integration scheme – the number of inversions of the original Wilson Dirac operator in the course of one trajectory – is about the same as the corresponding value in ref. . The values we have chosen for the mass parameters $`\mu _i`$ and the step numbers $`N_i`$ can be found in table 2 and one can see by comparing to ref. that the step numbers $`N_i`$ (or $`2N_i`$) are indeed quite similar.
The computation of the variation of $`S_\mathrm{G}`$ is, compared to the variations of the other action parts, almost negligible in terms of computer time. Therefore we set $`N_0`$ always large enough to ensure that the gauge part does not influence the acceptance rate negatively and we leave the gauge part out in the following discussion.
If one compares e.g. for simulation point $`C`$ the average norm of the fermionic forces, then one finds that it increases like $`1:40`$ ($`F_2:F_1`$). The maximal norm of the forces is accordingly strongly ordered, approximately like $`1:20`$. The corresponding relations in the step numbers we had to choose (see the values in table 2) increase only like $`1:6`$.
Therefore we conclude that the norm of the forces can indeed serve as a first criterion to tune the time scales and the preconditioning masses, by looking for a situation in which $`\mathrm{\Delta }\tau _iF_i_{\mathrm{max}}`$ is a constant independent of $`i`$. But, it cannot be the only criterion. Finally, the acceptance rate is determined by $`\mathrm{exp}(\mathrm{\Delta }H)`$, which depends in a more complicated way on the forces, see e.g. ref. .
It is well known that simulations with the HMC algorithm in particular for small quark masses become often unstable if the step sizes are too large. It is an important result that with the choice of parameters as can be found in table 2 our simulations appear to be very stable down to quark masses of the order of $`20\mathrm{MeV}`$. We did encounter only few large, but not exceptional, fluctuations in $`\mathrm{\Delta }H`$ during the runs. A typical history of $`\mathrm{\Delta }H`$ and the average plaquette value can be found in figure 3 for run $`C`$. Note that even a pion mass of about $`380\mathrm{MeV}`$ might be still to large to observe the asymptotic behavior of the algorithm.
All our runs reproduce the average plaquette expectation values quoted in ref. and, where available, in ref. within the statistical errors. Our results together with the number of measurements $`N_{\mathrm{meas}}`$ and the integrated autocorrelation time can be found in table 3. We also measured the values of the pseudo scalar, the vector and the current quark mass and our numbers agree within errors with the values quoted in refs. . These measurements were done on $`100`$ configurations separated by $`5`$ trajectories at each simulation point and we computed the aforementioned quantities with the methods explained in ref. . In order to improve the signal we used Jacobi smearing and random sources. Our results in physical units can be found in table 1. Note that the value for $`m_\mathrm{V}`$ at simulation point $`C`$ has to be taken with some caution, because the lattice time extend was a bit too small to be totally sure about the plateau.
In order to set the scale we determined the Sommer parameter $`r_0/a`$ . Our calculation used a static action with improved signal to noise ratio<sup>1</sup><sup>1</sup>1First results applying an improved static action in the computation of the static potential already appeared in . , the tree-level improved force and potential and we enhanced the overlap with the ground state of the potential using APE smeared spatial gauge links. The results can be found in table 1. For run $`A`$ and $`B`$ our values for $`r_0/a`$ agree very well within the errors with the value quoted in ref. . One should keep in mind, however, that the values for $`r_0/a`$ are computed on rather low statistics.
### 4.5 Algorithm performance
Any statement about the algorithm performance has to include autocorrelation times. Since different observables can have in general rather different autocorrelation times, also the algorithm performance is observable dependent. However, in the following we will use the plaquette integrated autocorrelation time to determine the performance. Note that other physical quantities such as hadron masses show in general very different autocorrelation times.
The values we measured for the plaquette integrated autocorrelation times can be found in table 3. It is interesting to observe that for runs $`A`$ and $`B`$ the values for the plaquette integrated autocorrelation times are smaller than the one found for the domain decomposition method. An explanation for this may be that in the algorithm of ref. a subset of all link variables is kept fixed during the molecular dynamics evolution, while in our HMC variant all link variables are updated.
Our value for $`\tau _{\mathrm{int}}(P)`$ for run $`A`$ is almost identical to the corresponding one found in ref. . In contrast, for simulation point $`B`$ our value is a factor of three smaller, which is only partly due to the significantly smaller acceptance rate of about $`60\%`$ quoted in ref. for this point.
A measure for the performance of the pure algorithm, implementation and machine independent, but incorporating the autocorrelation times is provided by the cost figure
$$\nu =10^3(2n+3)\tau _{\mathrm{int}}(P)$$
(4-5)
that has been introduced in ref. . $`n`$ in eq. (4-5) stands for either $`N_{N_{\mathrm{PF}}}`$ in case a LF integration scheme is used or $`2N_{N_{\mathrm{PF}}}`$ in case a SW integration scheme is used. $`\nu `$ represents the average number of inversions of the Wilson-Dirac operator with the physical mass in units of thousands as needed to generate a statistically independent value of the average plaquette. Hence, in giving values for $`\nu `$, we neglect the overhead coming from the remaining parts of the Hamiltonian.
Our values for $`\nu `$ together with the corresponding numbers from ref. and ref. are given in table 4. Compared to ref. our values for $`\nu `$ are smaller for simulation points $`A`$ and $`B`$ and comparable for run $`C`$. In contrast, the cost figure for the HMC algorithm with plain leap frog integration scheme is at least a factor $`10`$ larger than the values found for our HMC algorithm variant. This gain is, of course, what we aimed for by combining multiple time scale integration with mass preconditioning and hence confirms our expectation. Unfortunately, due to the large statistical uncertainties of the $`\nu `$ values it is not possible to give a scaling of the cost figure with the mass. This holds for our values of $`\nu `$ as well as the ones of ref. .
### 4.6 Simulation cost
Although the value of $`\nu `$ is a sensible performance measure for the algorithm itself, since it is independent of the machine, the actual implementation and the solver, it cannot serve to estimate the actual computer resources (costs) needed to generate one independent configuration. Assuming that the dominant contribution to the total cost stems from the matrix vector (MV) multiplications, we give in table 5 the average number of MV multiplications $`N_{\mathrm{MV}}`$ needed for the different pseudo fermion actions to evolve the system for one trajectory of length $`\tau =0.5`$. In addition we give the sum of these MV multiplications multiplied with the plaquette autocorrelation time together with the corresponding number from ref. .
In order to compare to the numbers of ref. we remark that the lattice time extent is $`T=40`$ in ref. compared to $`T=32`$ in our case, but we do not expect a large influence on the MV multiplications coming from this small difference. Large influence on the MV multiplications, however, we expect from ll-SSOR preconditioning that was used in ref. in combination with a chronological solver guess (CSG) .
Initially, when one compares the values of the cost figure for our HMC algorithm with the one of the plain leap frog algorithm as used in ref. , one might expect that the number of MV multiplications shows a similar behavior as a function of the quark mass. However, inspecting table 5, we see that in terms of MV multiplications at simulation point $`A`$ the HMC algorithm of ref. is only a factor of $`2`$ slower than the variant presented in this paper, while the values of $`\nu `$ are by a factor of about $`20`$ different. The reason for this is two-fold: On the one hand ll-SSOR preconditioning together with a CSG method is expected to perform better than only even/odd preconditioning. On the other hand we think that the quark mass at this simulation point is still not small enough to gain significantly from multiple time scale integration. This illustrates that indeed the value of $`\nu `$ is not immediately conclusive for the actual cost of the algorithm.
At simulation point $`B`$ the relative factor between the MV multiplications needed by the two algorithms is already about $`7`$. And finally, it is remarkable that for simulation point $`C`$ the costs with our HMC variant are still a factor of $`2`$ smaller than the costs for simulation point $`B`$ with the algorithm used in ref. , even though the masses are very different.
From this comparison we conclude that especially in the regime of small quark masses the HMC algorithm presented in this paper is significantly faster than a HMC algorithm with single time scale leap frog integration scheme.
By looking at table 5 one notices that especially for simulation point $`C`$ the number of MV multiplications needed for preconditioning is larger than the one needed for the physical operator. This comes from the fact that with the choice of algorithm parameters we have used the number of molecular dynamics steps for the mass preconditioned operator is large. This possibly indicates potential to further impove the performance by tuning the preconditioning masses and time scales.
We stress here again that the number of matrix vector operations is highly solver dependent, and therefore, every improvement to reduce the solver iterations will decrease the cost for one trajectory. Two promising improvements are the following:
* The use of a chronological inversion method :
The idea of the chronological inversion method (or similar methods ) is to optimize the initial guess for the solution used in the solver. To this end the history of $`N_{\mathrm{CSG}}`$ last solutions of the equation $`M^2\chi =\varphi `$ is saved and then a linear combination of the fields $`\chi _i`$ with coefficients $`c_i`$ is used as an initial guess for the next inversion. $`M`$ stands for the operator to be inverted and has to be replaced by the different ratios of operators used in this paper.
The coefficients $`c_i`$ are determined by solving
$$\underset{i}{}\chi _j^{}M^2\chi _ic_i=\chi _j^{}\varphi $$
(4-6)
with respect to the coefficients $`c_i`$. This is equivalent to minimizing the functional that is minimized by the CG inverter itself.
In ref. it was reported that with a chronological solver guess the number of MV multiplications can be reduced by a factor $`5`$ or even more. The gain is larger the smaller the size of the time steps is. But at the same time the reversibility violations increase at equal stopping criteria in the solver.
We have implemented the CSG method and tested its potential in the runs for this paper. On the one hand we see a significant reduction of MV multiplications on the small time scales, while the improvement for the large time scales is small, as expected.
On the other hand we observe that the reversibility violations increase significantly by one or two orders of magnitude in the Hamiltonian when the CSG is switched on and all other parameters are kept fixed. Therefore one has to adjust the residues in the solvers, which increases the number of MV multiplications again.
In total we found not more than a $`20\%`$ gain in matrix vector operations when a CSG is used. The largest gain is seen for the largest value of $`\kappa `$ under investigation. It is expected that this gain increases when the value of the bare physical mass is further reduced, because probably the size of the time steps must be further decreased.
* A different solver than the CG iterative solver, e.g. a solver using a Schwarz method as presented in ref. can also reduce the iteration numbers significantly. The method introduced in ref. is expected to be particularly useful for inverting the original fermion matrix with a small mass.
Finally, it is interesting to compare the number of matrix vector multiplications reported in table 5 with a HMC algorithm where mass preconditioning and multiple time scale improvements are switched off and CSG is not used. For instance for a simulation with a Sexton-Weingarten improved integration scheme at $`\kappa =0.15825`$ there are $`120`$ molecular dynamics steps needed to get acceptance. This corresponds to $`240`$ inversions of $`Q^2`$, which amounts to about $`720000`$ matrix vector multiplications. Compared to run $`C`$ this is at least a factor $`10`$ more. We did only a few trajectories to get an estimate for this number, so we cannot say anything about autocorrelation time.
Of course it would be interesting to compare also to a HMC algorithm with mass preconditioning but without multiple time scale integration. This, however, needs again a tuning of the mass parameters and would therefore be quite costly and we did not attempt to test this situation here.
### 4.7 Scaling with the mass
An important property of an algorithm for lattice QCD is the scaling of the costs with the simulated quark mass. The naive expectation is that the number of solver iterations grows like $`m_q^1`$ and also the number of molecular dynamics steps is proportional to $`m_q^1`$, see for instance ref. or ref. . Since also the integrated autocorrelation time is assumed to grow like $`m_q^1`$, it is expected that the HMC algorithm costs scale with the quark mass as $`m_q^3`$ or equivalently as $`m_{\mathrm{PS}}^6`$. In contrast, for our HMC algorithm variant we expect a much weaker scaling of $`\mathrm{\Delta }\tau `$ and also of the number of solver iterations. Indeed, we see that the costs for our HMC algorithm variant is consistent with a $`m_q^2`$ or $`m_{\mathrm{PS}}^4`$ behaviour when the autocorrelation time is taken into account.
We have translated the number of matrix vector multiplications from table 5 into costs in units $`\mathrm{Tflops}\mathrm{years}`$ and plotted the computer resources needed to generate $`1000`$ independent configurations of size $`24^3\times 40`$ at a lattice spacing of $`0.08\mathrm{fm}`$ as a function of $`m_{\mathrm{PS}}/m_\mathrm{V}`$ in figure 4(a) together with the results of ref. . Note that we have scaled our costs like $`(40/32)^{1.25}`$ corresponding to the expected volume dependence (cf. ) to match the different time extents and, moreover, we used the plaquette autocorrelation time as an estimate for the autocorrelation time.
The solid (dashed) line is not a fit to the data, but a function proportional to $`(m_{\mathrm{PS}}/m_\mathrm{V})^4`$ ($`(m_{\mathrm{PS}}/m_\mathrm{V})^6`$) with a coefficient that is fixed by the data point corresponding to the lightest pseudo scalar mass. This functional dependencies on $`(m_{\mathrm{PS}}/m_\mathrm{V})`$ describes the data reasonably well. However, from figure 4(a) it is not possible to decide on the value of the exponent in the quark mass dependence of the costs. But, it is clear from the figure that with multiple time scale integration and mass preconditioning the “wall” – which renders simulations at some point infeasible – is moved towards smaller values of the quark mass.
On a larger scale we can compare the extrapolations of our cost data to the formula given in ref.
$$C=K\left(\frac{m_{\mathrm{PS}}}{m_\mathrm{V}}\right)^{z_\pi }L^{z_L}a^{z_a},$$
(4-7)
where the constant $`K`$ can be found in ref. and $`z_\pi =6`$, $`z_L=5`$ and $`z_a=7`$. The result of this comparison is plotted in figure 4(b), which is an update of the “Berlin Wall” figure that can be found in ref. . We plot the simulation costs in units of $`\mathrm{Tflops}\mathrm{years}`$ versus $`m_{\mathrm{PS}}/m_\mathrm{V}`$, where we again scaled the numbers in order to match a lattice time extend of $`T=40`$. The dashed and the dotted lines are extrapolations from our data proportional to $`(m_{\mathrm{PS}}/m_\mathrm{V})^4`$ and $`(m_{\mathrm{PS}}/m_\mathrm{V})^6`$, respectively, again matching the data point corresponding to the lightest pseudo scalar mass. The solid line corresponds to eq. (4-7). In addition we plot data from staggered simulations as were used for the plot in ref. . That the corresponding points lie nearly on top of the dotted line is accidental.
Conservatively one can conclude from figure 4(b) that with the HMC algorithm described in this paper at least simulations with $`m_{\mathrm{PS}}/m_\mathrm{V}0.3`$ are feasible, even though $`L=1.93\mathrm{fm}`$ is too small for such values of the masses. Taking the more optimistic point of view by assuming that the costs scale with $`z_\pi =4`$, even simulation with the physical $`m_{\mathrm{PS}}/m_\mathrm{V}`$ ratio and a lattice spacing of $`0.08\mathrm{fm}`$ become accessible, with again the caveat that $`L/a`$ needs to be increased.
Independent of the value for $`z_\pi `$, figure 4(b) reveals that the costs for simulations with staggered fermions and with Wilson fermions in a comparable physical situation are of the same order of magnitude, if for the simulations with Wilson fermions an algorithm like the one presented in this work is used. It would be interesting to see whether the techniques applied in this paper work similarly well for staggered fermions.
We would like to point out that we did not try to tune the parameters to their optimal values. The aim of this paper was to give a first comparison of mass preconditioned HMC algorithm with multiple time scale integration to existing performance data, i.e. data for a HMC algorithm preconditioned by domain decomposition and data for the HMC algorithm variant of ref. . We are confident that there are still improvements possible by further tuning of the parameters in our variant of the HMC algorithm.
## 5 Conclusion
In this paper we have presented and tested a variant of the HMC algorithm combining multiple time scale integration with mass preconditioning (Hasenbusch acceleration). The aim of this paper was to perform a first investigation of the performance properties of this HMC algorithm by comparing it to other state of the art HMC algorithm variants at the same situation, i.e. for bare quark masses in the range of $`20`$ to $`60\mathrm{MeV}`$, a lattice spacing of about $`0.08\mathrm{fm}`$ and a lattice size of $`L2`$fm with two flavors of mass degenerate Wilson fermions.
We computed at each simulation point the expectation values of the plaquette and of the pion, the vector and the current quark masses finding full agreement with results in the literature. In order to set the scale we computed the Sommer scale , providing a value for $`r_0/a`$ also at the lowest quark mass we simulated and which has not been available in the literature so far.
We have shown that the additional mass parameters introduced for mass preconditioning can be arranged such that the force contributions from the different parts in the Hamiltonian are strictly ordered with respect to the absolute value of the force and that the most expensive part has the smallest contribution to the total force.
Using this result, it is possible to tune the time scales such that the performance of our variant in terms of the cost figure in eq. (4-5) is compatible to the one observed for the HMC algorithm with multiple time scales and domain decomposition as preconditioner introduced in ref. and clearly superior to the one for the HMC algorithm with a simple leap frog integration scheme as used in ref. .
While the cost figure provides a clean algorithm performance measure we also compare the simulation costs in units of $`\mathrm{Tflops}\mathrm{years}`$ to existing data. This comparison is summarized in an update of the “Berlin wall” plot of ref. , which can be found in figure 4. We could show that with the HMC algorithm presented in this paper the wall is moved towards smaller values of the quark mass and that simulations with a ratio of $`m_{\mathrm{PS}}/m_\mathrm{V}0.3`$ become feasible at a lattice spacing of around $`0.08\mathrm{fm}`$ and $`L2`$fm.
The HMC variant presented here has the advantage of being applicable to a wide variety of Dirac operators, including in principle also the overlap operator. In addition its implementation is straightforward, in particular in an already existing HMC code. We remark that the paralllelization properties of our HMC variant and the one of the algorithm presented in can be very different depending on whether a fine- or a coarse-grained massively parallel computer architecture is used.
From a stability point of view our results reveal that even for Wilson fermions it is very well possible to simulate quark masses of the order of $`20\mathrm{MeV}`$ when using the algorithmic ideas presented in this paper. We are presently simulating even smaller quark masses without practical problems, but the statistics is not yet adequate to say something definite.
The results presented in this paper are mostly based on empirical observations and on simulations for only one value of the coupling constant $`\beta =5.6`$. It remains to be seen how our HMC variant behaves for larger values of $`\beta `$, which, as well as smaller quark masses and theoretical considerations about the scaling properties with the quark mass needs further investigations.
Moreover, a more systematic study of the interplay between integration schemes, step sizes, (preconditioning and physical) masses and the simulation costs is needed. Those investigations will hopefully also provide a better understanding of the algorithm itself and its dynamics.
Finally, we think that there are further improvements possible by the usage of a Polynomial HMC (PHMC) algorithm . With such an algorithm one could treat the lowest eigenvalues of the Dirac operator exactly and/or by reweighting. In this set-up the large fluctuations in the force might be significantly reduced, if the lowest eigenvalues are responsible for those. Then it might be possible to further reduce the number of inversions of the badly conditioned physical operator needed to evolve the system.
## Acknowledgements
We thank M. Lüscher, I. Montvay and I. Wetzorke for helpful comments and discussions, the qq+q collaboration and in particular F. Farchioni, I. Montvay and E.E. Scholz for providing us their analysis program for the masses. We thank C. Destri and R. Frezzotti for giving us access to a PC cluster in Milano, where parts of the computations for this paper have been performed. We also thank the computer-centers at HLRN and at DESY Zeuthen for granting the necessary computer-resources, and M. Hasenbusch for leaving us his Wilson HMC code as a starting point. This work was supported by the DFG Sonderforschungsbereich/Transregio SFB/TR9-03. |
warning/0506/quant-ph0506014.html | ar5iv | text | # 1 Introduction
## 1 Introduction
A lot of developments and applications of the classical approaches of Gelfand, Levitan and Marchenko for the solution of the inverse-scattering problem at fixed angular momentum exist, and we have several of excellent reviews on the subject . The direct application of these approaches to construction of local two-body potentials which are phase equivalent to the effective potentials occurring in theories describing reactions of composite particles is impossible. For, such potential must be complex in order to reproduce the loss of flux above the inelastic threshold. But it must reproduce the real phase shifts below this threshold and must be real itself. These requirements are incompatible for potentials being energy independent by construction. For very low threshold these approaches are applicable and produce energy independent complex potentials . In the general case the empirical energy dependent optical potentials are usually inferred by fitting of the parameters of an assumed analytic potential . This approach has two major shortcomings: a complexity and inconvenience of fitting simultaneously many nonlinear parameters; and lack of correlation of the parameters obtained at various energies.
In this paper we develop an inversion method that is free of these shortcomings. The method is based on a fixed-$`l`$ inverse scattering theory and on a special parameterization of the optical potential. The proposed procedure is a two-step process. In the first step the phase shift data are used to determine a real potential via solution of the Marchenko equation. At this point we use a diagonal Padé approximant $`[M/M]`$ of the corresponding unitary $`S`$-matrix. In the second step the imaginary part of the potential is determined via the phase equation of the variable phase approach . We assume that the real and the imaginary parts of the optical potential are proportional. The value of the proportionality coefficient is predicted by the phase equation and is refined by the iterative algorithm. We develop this method for single and for coupled partial waves. The whole procedure is applied to analyze $`{}_{}{}^{1}S_{0}^{}`$ $`NN`$, $`{}_{}{}^{3}SD_1`$ $`NN`$ data and to analyze $`P31`$ $`\pi N`$ and $`S01`$ $`K+N`$ data. These analyses demonstrate that prediction for the proportionality coefficient from the phase equation is very close to a precise value that reproduce the experimental loss of flux.
The plan of the paper is as follows. In Sect. 2 we describe the inverse scattering techniques based on the Marchenko integral equation. The used diagonal Padé approximants of the corresponding $`S`$-matrix allow an analytical solution of the Marchenko integral equation . For single partial wave the general solution was presented in . We present a solution for coupled partial waves. These techniques produce real local potentials from phase shift analysis data. In Sect. 3 we consider the phase equation. We investigate how the $`S`$-matrix is changing with certain change of the potential. This consideration shows advantages of proportionality of the real and imaginary parts of the optical potential. In Sect. 4 the feasibility of the method is shown in the examples of analyses of $`NN`$, $`\pi ^{}N`$ and $`K^+N`$ scattering data.
## 2 Inversion algorithm
The Marchenko inverse scattering theory is viewed in detail in Refs. . We shall, therefore, only briefly describe this formalism.
The input data of the Marchenko inversion are
$$\left\{S\left(q\right),\left(0<q<+\mathrm{}\right),\text{ }\stackrel{~}{q}_j,\text{ }M_j,\text{ }j=1,\mathrm{},n_\text{b}\right\},$$
(1)
where $`S\left(q\right)`$ is the scattering matrix dependant on the relative momentum $`q`$, $`q^2=Em`$, $`\stackrel{~}{q}_j^2=mE_j0,E_j`$ is the energy of the j-th bound state, so that $`ı\stackrel{~}{q}_j0`$, $`m`$ is the reduced mass. The $`M_j`$ matrices give the asymptotic behavior of the corresponding normalized bound states.
We proceed from the Marchenko equation for single channel
$$F(x,y)+L(x,y)+\underset{x}{\overset{+\mathrm{}}{}}L(x,t)F(t,y)𝑑t=0,$$
(2)
where the input kernel is given by
$$F(x,y)=\frac{1}{2\pi }\underset{\mathrm{}}{\overset{+\mathrm{}}{}}h_l^+\left(qx\right)\left(IS\left(q\right)\right)h_l^+\left(qy\right)𝑑q+\underset{j=1}{\overset{n_\text{b}}{}}M_j^2h_l^+\left(iq_jx\right)h_j^+\left(iq_jy\right),$$
(3)
$`h_l^+\left(z\right)`$ are the Riccati-Hankel functions.
The output kernel $`L(x,y)`$ gives the reconstructed potential
$$V\left(r\right)=\frac{dL(r,r)}{dr}.$$
(4)
This local energy independent operator $`V\left(r\right)`$ links the Marchenko equation (2) and the radial Schrödinger equation of a fixed angular momentum,
$$\left[\frac{d^2}{dr^2}+\frac{l(l+1)}{r^2}+V(r)\right]\psi (r,q)=q^2\psi (r,q).$$
(5)
The scattering matrix $`S\left(q\right)`$, matrices $`M_j`$ and energies $`E_j`$ are the output data of the direct scattering problem associated with the Schrödinger equation (5).
It has been known for several decades now that $`S`$ matrices rational in $`q`$ (ratio of polynomials) correspond to potentials known as Bargmann potentials expressible in terms of the elementary functions . Such fraction may have the same truncated Taylor series as the $`S`$ matrix it represents. It is then called a Padé approximant. Conjectures and theorems concerning the convergence and analytic continuation properties of Padé approximants are collected in . For single partial wave the general solution of the Marchenko equation via Padé approximant of the $`S`$-matrix was presented in and in . We shall, therefore, only present it briefly and turn to the case of coupled partial waves.
A diagonal Padé approximant of the $`S`$-matrix is given by
$$S\left(q\right)=e^{2ı\delta }=\frac{f_2\left(q\right)ıf_1\left(q\right)}{f_2\left(q\right)+ıf_1\left(q\right)}$$
(6)
$`f_1\left(q\right)`$ and $`f_2\left(q\right)`$ are an odd and even polynomials of $`q`$, which do not turn to zero at the real axis simultaneously.
This approximant leads to the following expression for the phase shifts $`\delta \left(q\right)`$
$$tg\left(\delta \left(q\right)\right)=\frac{f_1\left(q\right)}{f_2\left(q\right)}$$
(7)
Inasmuch as $`lim_q\mathrm{}\delta (q)1/q`$ for regular potentials, it is evident that degree of $`f_1(q)`$ must be less than degree of $`f_2(q)`$ by 1. Let us select $`N`$ discrete momenta $`q_i`$ such that the corresponding $`\delta \left(q_i\right)=\delta _i`$ are known. Use of these values in eq. (7) transforms the latter into a set of inhomogeneous linear equations from which $`N`$ coefficients of the polynomials $`f_1\left(q\right)`$ and $`f_2\left(q\right)`$ can be determined. This is a usual strategy , but since any set of $`q_i`$ is experimentally limited from above ($`q_i<q_{max}`$) there is some uncertainty in determination of $`S(q)`$. Even if a set of $`q_i`$ is dense and the agrement between the data and the used approximant is excellent the arbitrary behavior of $`S(q)`$ above $`q_{max}`$ guarantees that a solution of the inverse problem is arbitrary as well. We assume that about and above $`q_{max}`$ the true $`\delta (q)`$ depends only slightly on details of the potential and does not depend on its asymptotic at $`r>r_{max}`$. Then we take a model potential (in our calculations $`V_{model}(r)=A\mathrm{exp}(br)`$) and fit parameters ($`A`$ and $`b`$) so that
$$\delta _{model}(q_{max})\delta (q_{max}),$$
(8)
here signs $``$ mean that about $`q=q_{max}`$ the chosen model potential gives a phase curve which goes inside error bars. This means that we take $`\delta _{model}(q)`$ as an asymptotic for $`\delta (q)`$ when $`q>q_{max}`$. In the segment $`[0,q_{max}]`$ the $`\delta (q)`$ is defined by some spline that approximates the data points $`\{q_i,\delta _i\}`$. In this way we may control the Padé fit in the line segment $`[0,Q_{max}]`$, where $`Q_{max}`$ is arbitrary and is fixed by convergence of the whole inversion procedure. The needed accuracy of approximant is attained by increasing of $`N`$ which in turn defines degrees of $`f_1(q)`$ and $`f_2(q)`$.
Approximant (7) leads to a degenerate input kernel $`F(x,y)`$. We calculate the integral in eq. (3) using the residue theorem. For approximant (6) the result of the integration is
$$\begin{array}{c}F(x,y)=ı\underset{i=1}{\overset{n_{\text{pos}}}{}}Res\left[h_l^+\left(qx\right)\left(IS\left(q\right)\right)h_l^+\left(qy\right)\right]|_{q=\beta _i}+\underset{i=1}{\overset{n_\text{b}}{}}M_i^2h_l^+\left(\stackrel{~}{q}_ix\right)h_l^+\left(\stackrel{~}{q}_iy\right)=\\ =\underset{i=1}{\overset{n_{\text{pos}}}{}}b_ih_l^+\left(\beta _ix\right)h_l^+\left(\beta _iy\right)+\underset{i=1}{\overset{n_\text{b}}{}}M_i^2h_l^+\left(\stackrel{~}{q}_ix\right)h_l^+\left(\stackrel{~}{q}_iy\right)=\underset{j=1}{\overset{n}{}}b_jh_l^+\left(\beta _jx\right)h_l^+\left(\beta _jy\right),\end{array}$$
(9)
where $`\beta _i`$ ($`i=1,\mathrm{},n_{\text{pos}})`$ are all $`S`$-matrix poles with $`\mathrm{}\beta _i>0`$, $`\beta =\{\beta _1,..\beta _{n_{pos}},\stackrel{~}{q_1},\mathrm{}\stackrel{~}{q_{n_b}}\}`$, $`n=n_{pos}+n_b`$. We assume that all poles are of first order so that
$$\begin{array}{c}Res\left[h_l^+\left(qx\right)\left(IS\left(q\right)\right)h_l^+\left(qy\right)\right]|_{q=\beta _i}=2ıRes\left[h_l^+\left(qx\right)\frac{f_1\left(q\right)}{f_2\left(q\right)+ıf_1\left(q\right)}h_l^+\left(qy\right)\right]|_{q=\beta _i}=\\ =2ı\frac{f_1\left(\beta _i\right)}{f_2^{}\left(\beta _i\right)+ıf_1^{}\left(\beta _i\right)}h_l^+\left(\beta _ix\right)h_l^+\left(\beta _iy\right)=b_ih_l^+\left(\beta _ix\right)h_l^+\left(\beta _iy\right),\end{array}$$
(10)
here we have denoted $`f_i^{}(q)=df_i(q)/dq`$, ($`i=1,\mathrm{\hspace{0.17em}2}`$).
In this case the input kernel of eq. (2) is a degenerate one as well as its output kernel
$$L(x,y)=\underset{i=1}{\overset{n}{}}P_i\left(x\right)h_l^+\left(\beta _iy\right),$$
(11)
where $`P_i\left(x\right)`$ are unknown coefficients. Substitution of (9) and (11) into (2) yields
$$\underset{i=1}{\overset{n}{}}h_l^+\left(\beta _iy\right)\left(b_ih_l^+\left(\beta _ix\right)+P_i(x)+b_i\underset{k=1}{\overset{n}{}}P_k\left(x\right)\underset{x}{\overset{+\mathrm{}}{}}h_l^+\left(\beta _kt\right)h_l^+\left(\beta _it\right)𝑑t\right)=0.$$
(12)
Linear independence of the $`h_l^+\left(\beta _iy\right)`$ implies that
$$b_ih_l^+\left(\beta _ix\right)+P_i(x)+b_i\underset{k=1}{\overset{n}{}}P_k\left(x\right)\underset{x}{\overset{+\mathrm{}}{}}h_l^+\left(\beta _kt\right)h_l^+\left(\beta _it\right)𝑑t=0,$$
(13)
or
$$\underset{k=1}{\overset{n}{}}A_{ik}\left(x\right)P_k\left(x\right)=D_i\left(x\right)(i=1,..,n),$$
(14)
where $`D_i\left(x\right)=b_ih_l^+\left(\beta _ix\right)`$ and after applying Riccati-Hankel integration formulas (see Appendix) in (13) we have
$$\begin{array}{cc}A_{ik}=& \{\begin{array}{c}1+b_ix(\left(h_l^+\left(\beta _ix\right)\right)^2h_{l1}^+\left(\beta _ix\right)h_{l+1}^+\left(\beta _ix\right))/2,fori=k\\ b_i\frac{\beta _ih_{l1}^+\left(\beta _ix\right)h_l^+\left(\beta _kx\right)\beta _kh_l^+\left(\beta _ix\right)h_{l1}^+\left(\beta _kx\right)}{\beta _i^2\beta _k^2},forik.\end{array}\end{array}$$
(15)
The functional coefficients $`P_i\left(x\right)`$ are defined by (14)
$$P_i\left(x\right)=\left(A^1D\right)_i.$$
(16)
Finally we derive $`L(x,y)`$ and the potential $`V\left(r\right)`$ from (11) and (4).
In case of two coupled channels we present only sketchy derivations because of their awkwardness. In this case the system of the partial Schrödinger equations is
$$\left(\frac{d^2}{dr^2}+V\left(r\right)+\left(\begin{array}{cccccccccccccccccccc}\frac{l_1\left(l_1+1\right)}{r^2}& 0& & & & & & & & & & & & & & & & & & \\ 0& \frac{l_2\left(l_2+1\right)}{r^2}& & & & & & & & & & & & & & & & & & \end{array}\right)\right)\left(\begin{array}{cccccccccccccccccccc}\chi _1(r)& & & & & & & & & & & & & & & & & & & \\ \chi _2(r)& & & & & & & & & & & & & & & & & & & \end{array}\right)=q^2\left(\begin{array}{cccccccccccccccccccc}\chi _1(r)& & & & & & & & & & & & & & & & & & & \\ \chi _2(r)& & & & & & & & & & & & & & & & & & & \end{array}\right),$$
(17)
$$V\left(r\right)=\left(\begin{array}{cccccccccccccccccccc}V_1\left(r\right)& V_T\left(r\right)& & & & & & & & & & & & & & & & & & \\ V_T\left(r\right)& V_2\left(r\right)& & & & & & & & & & & & & & & & & & \end{array}\right),$$
(18)
where $`V_1\left(r\right)`$, $`V_2\left(r\right)`$ are potentials in channels 1 and 2, $`V_T\left(r\right)`$ is potential coupling them, $`\chi _1(r)`$ and $`\chi _2(r)`$ are channel wave functions.
By analogy with (6) we approximate the $`S`$-matrix by the following expression
$$\begin{array}{c}S(x)=\left(\begin{array}{cccccccccccccccccccc}\mathrm{exp}\left(2ı\delta _1\right)\mathrm{cos}2\epsilon & i\mathrm{exp}\left(ı\left(\delta _1+\delta _2\right)\right)\mathrm{sin}2\epsilon & & & & & & & & & & & & & & & & & & \\ ı\mathrm{exp}\left(ı\left(\delta _1+\delta _2\right)\right)\mathrm{sin}2\epsilon & \mathrm{exp}\left(2ı\delta _2\right)\mathrm{cos}2\epsilon & & & & & & & & & & & & & & & & & & \end{array}\right)=\hfill \\ \\ =\left(\begin{array}{cccccccccccccccccccc}\left(\frac{f_2^{\left(1\right)}\left(q\right)ıf_1^{\left(1\right)}\left(q\right)}{f_2^{\left(1\right)}\left(q\right)+ıf_1^{\left(1\right)}\left(q\right)}\right)^2\frac{\left(f_2^{\left(12\right)}\left(q\right)\right)^2\left(f_1^{\left(12\right)}\left(q\right)\right)^2}{\left(f_2^{\left(12\right)}\left(q\right)\right)^2+\left(f_1^{\left(12\right)}\left(q\right)\right)^2}& 2ı\frac{f_2^{\left(12\right)}\left(x\right)f_1^{\left(12\right)}\left(x\right)}{\left(f_2^{\left(12\right)}\left(q\right)\right)^2+\left(f_1^{\left(12\right)}\left(q\right)\right)^2}\underset{j=1,2}{}\frac{f_2^{\left(j\right)}\left(q\right)ıf_1^{\left(j\right)}\left(q\right)}{f_2^{\left(j\right)}\left(q\right)+ıf_1^{\left(j\right)}\left(q\right)}& & & & & & & & & & & & & & & & & & \\ 2i\frac{f_2^{\left(12\right)}\left(q\right)f_1^{\left(12\right)}\left(q\right)}{\left(f_2^{\left(12\right)}\left(q\right)\right)^2+\left(f_1^{\left(12\right)}\left(q\right)\right)^2}\underset{j=1,2}{}\frac{f_2^{\left(j\right)}\left(q\right)if_1^{\left(j\right)}\left(q\right)}{f_2^{\left(j\right)}\left(q\right)+if_1^{\left(j\right)}\left(q\right)}& \left(\frac{f_2^{\left(2\right)}\left(q\right)ıf_1^{\left(2\right)}\left(q\right)}{f_2^{\left(2\right)}\left(q\right)+ıf_1^{\left(2\right)}\left(q\right)}\right)^2\frac{\left(f_2^{\left(12\right)}\left(q\right)\right)^2\left(f_1^{\left(12\right)}\left(q\right)\right)^2}{\left(f_2^{\left(12\right)}\left(q\right)\right)^2+\left(f_1^{\left(12\right)}\left(q\right)\right)^2}& & & & & & & & & & & & & & & & & & \end{array}\right)\hfill \end{array}$$
(19)
This is again the most general Padé approximant for the $`S`$-matrix. It was used in in an other form, but corresponding analytical solution of the inverse scattering problem was not presented.
The coefficients of this Padé approximant are determined from the equations analogous to (20)
$$tg\left(\frac{\delta _i\left(q\right)}{2}\right)=\frac{f_1^{(i)}\left(q\right)}{f_2^{(i)}\left(q\right)},i=1,2$$
(20)
and
$$tg\left(\epsilon \left(q\right)\right)=\frac{f_1\left(q\right)}{f_2\left(q\right)}.$$
(21)
The generalized Marchenko equation for coupled channels formally has the former view
$$L(x,y)+F(x,y)+\underset{x}{\overset{+\mathrm{}}{}}L(x,t)F(t,y)𝑑t=0,$$
(22)
but functions involved are matrices $`\left(2\times 2\right)`$
$$F(x,y)=\frac{1}{2\pi }\underset{\mathrm{}}{\overset{+\mathrm{}}{}}H\left(qx\right)\left[IS\left(q\right)\right]H\left(qy\right)𝑑q+\underset{i=1}{\overset{n_b}{}}H\left(\beta _ix\right)M_iH\left(\beta _iy\right),$$
(23)
where
$$H\left(x\right)=\left(\begin{array}{cccccccccccccccccccc}h_{l_1}^+\left(x\right)& 0& & & & & & & & & & & & & & & & & & \\ 0& h_{l_2}^+\left(x\right)& & & & & & & & & & & & & & & & & & \end{array}\right),I=\left(\begin{array}{cccccccccccccccccccc}1& 0& & & & & & & & & & & & & & & & & & \\ 0& 1& & & & & & & & & & & & & & & & & & \end{array}\right).$$
(24)
Insertion of (19) into (23) and applying of the residue theorem yields
$$\begin{array}{c}F(x,y)=ı\underset{i=1}{\overset{n_{pos}}{}}Res\left[H\left(qx\right)\left(IS\left(q\right)\right)H\left(qy\right)\right]|_{q=\beta _i}+\underset{i=1}{\overset{n_\text{b}}{}}H\left(\beta _ix\right)M_i^2H\left(\beta _iy\right)=\hfill \\ =\underset{i=1}{\overset{n}{}}H\left(\beta _ix\right)Q_i^1H\left(\beta _iy\right)+\underset{i=1}{\overset{n_{pos}^{(2)}}{}}xH^{}\left(\beta _ix\right)Q_i^2H\left(\beta _iy\right)+\underset{i=1}{\overset{n_{pos}^{(2)}}{}}H\left(\beta _ix\right)Q_i^2H^{}\left(\beta _iy\right)y,\hfill \end{array}$$
(25)
where $`\beta _i`$ ($`i=1,\mathrm{},n_{\text{pos}})`$ are all $`S`$-matrix poles with $`\mathrm{}\beta _i>0`$, $`\beta _i`$ ($`i=1,\mathrm{},n_{\text{pos}}^{\left(2\right)})`$ are poles of the second order, $`\beta =\{\beta _1,..\beta _{n_{\text{pos}}^{(2)}},..\beta _{n_{pos}},\stackrel{~}{q_1},\mathrm{}\stackrel{~}{q_{n_b}}\}`$, $`n=n_{pos}+n_b`$,
$$H^{}\left(x\right)=\left(\begin{array}{cccccccccccccccccccc}dh_{l_1}^+\left(x\right)\mathrm{/}dx& 0& & & & & & & & & & & & & & & & & & \\ 0& dh_{l_2}^+\left(x\right)\mathrm{/}dx& & & & & & & & & & & & & & & & & & \end{array}\right).$$
We note that there are poles of the first as well as of the second order in the diagonal matrix elements and there are poles only of the first order in the off-diagonal matrix elements. Poles of the second order in the diagonal elements are poles of the first order in the off-diagonal matrix elements, and they must be enumerated twice. $`Q_i^j\left(j=1,2\right)`$ are constant matrices they are trivial but cumbersome therefore, we do not give them.
We solve eq. (22) using substitution
$$L(x,y)=\underset{i=1}{\overset{n}{}}P_i\left(x\right)H\left(\beta _iy\right)+\underset{i=1}{\overset{n}{}}N_i\left(x\right)yH^{}\left(\beta _iy\right),$$
(26)
where $`P_i\left(x\right)`$, $`N_i\left(x\right)`$ are unknown functional $`\left(2\times 2\right)`$ matrix-coefficients. Linear independence of the $`H\left(\beta _iy\right)`$ and $`yH^{}\left(\beta _iy\right)`$ implies that
$$\begin{array}{c}\hfill \underset{i}{}P_i\left(x\right)Q_{ij}^3\left(x\right)+\underset{i}{}N_i\left(x\right)Q_{ij}^5\left(x\right)=H\left(\beta _jx\right)Q_j^1+xH^{}\left(\beta _jx\right)Q_j^2\\ \hfill \underset{i}{}N_i\left(x\right)Q_{ij}^6\left(x\right)+\underset{i}{}P_i\left(x\right)Q_{ij}^4\left(x\right)=H\left(\beta _jx\right)Q_j^2\end{array}$$
(27)
where
$$Q_{ij}^3\left(x\right)=I\delta _{ij}+\underset{x}{\overset{+\mathrm{}}{}}H\left(\beta _it\right)H\left(\beta _jt\right)𝑑t\times Q_j^1+\underset{x}{\overset{+\mathrm{}}{}}tH\left(\beta _it\right)H^{}\left(\beta _jt\right)𝑑t\times Q_j^2$$
$$Q_{ij}^4\left(x\right)=\underset{x}{\overset{+\mathrm{}}{}}H\left(\beta _it\right)H\left(\beta _jt\right)𝑑t\times Q_j^2$$
(28)
$$Q_{ij}^5\left(x\right)=\underset{x}{\overset{+\mathrm{}}{}}tH^{}\left(\beta _it\right)H\left(\beta _jt\right)𝑑t\times Q_j^1+\underset{x}{\overset{+\mathrm{}}{}}t^2H^{}\left(\beta _it\right)H^{}\left(\beta _jt\right)𝑑t\times Q_j^2$$
$$Q_{ij}^6\left(x\right)=I\delta _{ij}+\underset{x}{\overset{+\mathrm{}}{}}tH^{}\left(\beta _it\right)H\left(\beta _jt\right)𝑑t\times Q_j^2,$$
Integrals of expressions (28) are presented in Appendix. Matrix equations (27) can be trivially reduced to scalar linear equations. Having solved this linear equation system we get the sought-for potential from (26) and (4).
The multichannel generalization is trivial.
## 3 The optical potential
In this section we consider changes of $`S`$-matrix that are induced by certain transformation of real potential.
First we consider the one channel problem.
The phase equation for the initial potential $`V^0\left(r\right)`$ obtained by some inversion procedure (from Marchenko equation in our calculations) is
$$\delta _l^{(0)}=\frac{1}{q}\underset{0}{\overset{\mathrm{}}{}}V^{(0)}\left(r\right)\widehat{D}_l^2\left(qr\right)\mathrm{sin}^2\left(\widehat{\delta }_l(qr)+\delta ^{(0)}\left(r\right)\right)𝑑r,$$
(29)
where $`\widehat{D}_l\left(z\right)`$ and $`\widehat{\delta }_l\left(z\right)`$ are Riccati-Bessel amplitude and phase correspondingly
$$\widehat{D}_l(x)=\sqrt{j_l^2(x)+n_l^2(x)},$$
(30)
$$\widehat{\delta }_l(x)=\mathrm{arctan}(j_l(x)/n_l(x))$$
(31)
Let us consider the complex-valued potential $`V^{(1)}\left(r\right)`$ obtained from $`V^0\left(r\right)`$ by transformation
$$V^{(1)}\left(r\right)=\left(1+i\alpha \right)V^{(0)}\left(r\right),$$
(32)
where $`\alpha `$ is some real parameter. Such parametrization was used in but without analysis ($`\alpha `$ was fitted). Evidently the phase equation for this potential is
$$\delta ^{(1)}=\frac{1}{q}\left(1+i\alpha \right)\underset{0}{\overset{\mathrm{}}{}}V^{(0)}\left(r\right)\widehat{D}_l^2\left(qr\right)\mathrm{sin}^2\left(\widehat{\delta }_l(qr)+\delta ^{(1)}\left(r\right)\right)𝑑r.$$
(33)
From eqs. (29) and (33) we get
$$\begin{array}{c}\delta ^{(1)}\left(1+i\alpha \right)\delta ^{(0)}=\hfill \\ =\frac{1+i\alpha }{q}\underset{0}{\overset{\mathrm{}}{}}V^{(0)}\left(r\right)\widehat{D}_l^2\left(qr\right)\left(\mathrm{sin}^2\left(\widehat{\delta }_l(qr)+\delta ^{(1)}\left(r\right)\right)\mathrm{sin}^2\left(\widehat{\delta }_l(qr)+\delta ^{(0)}\left(r\right)\right)\right)𝑑r=\hfill \\ =\frac{1+i\alpha }{q}\underset{0}{\overset{\mathrm{}}{}}V^{(0)}\left(r\right)\widehat{D}_l^2\left(qr\right)\underset{¯}{\mathrm{sin}\left(2\widehat{\delta }_l(qr)+\delta ^{(1)}\left(r\right)+\delta ^{(0)}\left(r\right)\right)}\mathrm{sin}\left(\delta ^{(1)}\left(r\right)\delta ^{(0)}\left(r\right)\right)𝑑r\hfill \end{array}$$
(34)
For smooth enough potentials the right side of eq. (34) rapidly decreases comparing with $`\delta ^{(0)}`$ and $`\delta ^{(1)}`$, because there is a rapidly oscillating around zero function under the integral in (34) (underlined). Its frequency behaves as $`2q`$ for big $`q`$ (see (31)). Then as the first approximation we may take
$$\delta ^{(1)}\left(1+i\alpha \right)\delta ^{(0)}=\delta _R+i\delta _I.$$
(35)
For inelastic scattering the $`S`$-matrix is expressed through the real inelasticity parameter $`\rho `$ and the real phase shift $`\delta `$
$$S=\mathrm{cos}^2\left(\rho \right)e^{2i\delta }=e^{2i\left(\delta _R+i\delta _I\right)},$$
(36)
so we easily arrive at
$$\delta _R=\delta \delta ^{(0)},$$
(37)
$$\mathrm{cos}^2\left(\rho \right)e^{2\alpha \delta ^{(0)}},$$
(38)
whence it follows that $`\alpha \delta 0`$. The formula (38) allows to calculate the parameter $`\alpha `$ from the known values $`\rho `$ and $`\delta ^{(0)}\delta `$.
Eqs. (34-35) imply that
$$\underset{0}{\overset{\mathrm{}}{}}V^{(0)}\left(r\right)\widehat{D}_l^2\left(qr\right)\mathrm{sin}^2\left(\widehat{\delta }_l(qr)+\delta ^{(0)}\left(r\right)\right)𝑑r\underset{0}{\overset{\mathrm{}}{}}V^{(0)}\left(r\right)\widehat{D}_l^2\left(qr\right)\mathrm{sin}^2\left(\widehat{\delta }_l(qr)+\delta ^{(1)}\left(r\right)\right)𝑑r,$$
(39)
Consideration of the coupled partial waves is more complicated. The initial real potential is
$$V^{(0)}\left(r\right)=\left(\begin{array}{cccccccccccccccccccc}V_1^{(0)}& V_T^{(0)}& & & & & & & & & & & & & & & & & & \\ V_T^{(0)}& V_2^{(0)}& & & & & & & & & & & & & & & & & & \end{array}\right).$$
(40)
The equations for eigenphases and mixing parameter of potential (40) are
$`\delta _1^{(1)}=I_{11}^{(0)}+I_{12}^{(0)}+I_{13}^{(0)}`$ (41)
$`\delta _2^{(1)}=I_{21}^{(0)}+I_{22}^{(0)}+I_{23}^{(0)}`$ (42)
$`ϵ^{(1)}=I_{31}^{(0)}+I_{32}^{(0)}+I_{33}^{(0)},`$ (43)
where
$`I_{11}^{(0)}={\displaystyle \frac{1}{q}}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}𝑑rV_1^{(0)}(r)\mathrm{cos}^2ϵ^{(0)}(r)\widehat{D}_{l_1}^2(qr)\mathrm{sin}^2(\widehat{\delta }_{l_1}(qr)+\delta _1^{(0)}(r))`$
$`I_{12}^{(0)}={\displaystyle \frac{1}{q}}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}𝑑rV_2^{(0)}(r)\mathrm{sin}^2ϵ^{(0)}(r)\widehat{D}_{l_2}^2(qr)\mathrm{sin}^2(\widehat{\delta }_{l_2}(qr)+\delta _1^{(0)}(r))`$
$`I_{13}^{(0)}={\displaystyle \frac{1}{q}}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}𝑑rV_T^{(0)}(r)\mathrm{sin}2ϵ^{(0)}(r)\widehat{D}_{l_2}(qr)\mathrm{sin}(\widehat{\delta }_{l_2}(qr)+\delta _1^{(0)}(r))\widehat{D}_{l_1}(qr)\mathrm{sin}(\widehat{\delta }_{l_1}(qr)+\delta _1^{(0)}(r))`$ (44)
$`I_{21}^{(0)}={\displaystyle \frac{1}{q}}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}𝑑rV_1^{(0)}(r)\mathrm{sin}^2ϵ^{(0)}(r)\widehat{D}_{l_1}^2(qr)\mathrm{sin}^2(\widehat{\delta }_{l_1}(qr)+\delta _2^{(0)}(r))`$
$`I_{22}^{(0)}={\displaystyle \frac{1}{q}}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}𝑑rV_2^{(0)}(r)\mathrm{cos}^2ϵ^{(0)}(r)\widehat{D}_{l_2}^2(qr)\mathrm{sin}^2(\widehat{\delta }_{l_2}(qr)+\delta _2^{(0)}(r))`$
$`I_{23}^{(0)}={\displaystyle \frac{1}{q}}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}𝑑rV_T^{(0)}(r)\mathrm{sin}2ϵ^{(0)}(r)\widehat{D}_{l_2}(qr)\mathrm{sin}(\widehat{\delta }_{l_2}(qr)+\delta _2^{(0)}(r))\widehat{D}_{l_1}(qr)\mathrm{sin}(\widehat{\delta }_{l_1}(qr)+\delta _2^{(0)}(r))`$ (45)
$`I_{31}^{(0)}={\displaystyle \frac{1}{2q}}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\mathrm{sin}2ϵ^{(0)}(r)dr}{\mathrm{sin}(\delta _1^{(0)}(r)\delta _2^{(0)}(r))}}V_1^{(0)}(r)\widehat{D}_{l_1}^2(qr)\mathrm{sin}(\widehat{\delta }_{l_1}(qr)+\delta _1^{(0)}(r))\mathrm{sin}(\widehat{\delta }_{l_1}(qr)+\delta _2^{(0)}(r))`$
$`I_{31}^{(0)}={\displaystyle \frac{1}{2q}}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\mathrm{sin}2ϵ^{(0)}(r)dr}{\mathrm{sin}(\delta _1^{(0)}(r)\delta _2^{(0)}(r))}}V_2^{(0)}(r)\widehat{D}_{l_2}^2(qr)\mathrm{sin}(\widehat{\delta }_{l_2}(qr)+\delta _1^{(0)}(r))\mathrm{sin}(\widehat{\delta }_{l_2}(qr)+\delta _1^{(0)}(r))`$
$`I_{31}^{(0)}={\displaystyle \frac{1}{2q}}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{V_T^{(0)}(r)\widehat{D}_{l_1}(qr)\widehat{D}_{l_2}(qr)dr}{\mathrm{sin}(\delta _1^{(0)}(r)\delta _2^{(0)}(r))}}[\mathrm{cos}2ϵ^{(0)}(r)\mathrm{sin}(\widehat{\delta }_{l_1}(qr)+\delta _1^{(0)}(r))\mathrm{sin}(\widehat{\delta }_{l_2}(qr)+\delta _2^{(0)}(r))`$
$`{\displaystyle \frac{1}{2}}(\mathrm{cos}2ϵ^{(0)}(r)1)\mathrm{sin}(\delta _1^{(0)}(r)\delta _2^{(0)}(r))\mathrm{sin}(\widehat{\delta }_{l_1}(qr)\widehat{\delta }_{l_2}(qr))]`$ (46)
By analogy with the one channel case the following generalization for the optical potential is derived
$$V^{(1)}\left(r\right)=\left(\begin{array}{cccccccccccccccccccc}\left(1+i\alpha _1\right)V_1^{(0)}& \left(1+i\alpha _3\right)V_T^{(0)}& & & & & & & & & & & & & & & & & & \\ \left(1+i\alpha _3\right)V_T^{(0)}& \left(1+i\alpha _2\right)V_2^{(0)}& & & & & & & & & & & & & & & & & & \end{array}\right).$$
(47)
Evidently the phase equations for this potential is
$`\delta _1^{(1)}=\left(1+i\alpha _1\right)I_{11}^{(1)}+\left(1+i\alpha _2\right)I_{12}^{(1)}+\left(1+i\alpha _3\right)I_{13}^{(1)}`$ (48)
$`\delta _2^{(1)}=\left(1+i\alpha _1\right)I_{21}^{(1)}+\left(1+i\alpha _2\right)I_{22}^{(1)}+\left(1+i\alpha _3\right)I_{23}^{(1)}`$ (49)
$`ϵ^{(1)}=\left(1+i\alpha _1\right)I_{31}^{(1)}+\left(1+i\alpha _2\right)I_{32}^{(1)}+\left(1+i\alpha _3\right)I_{33}^{(1)}`$ (50)
Integrals $`I_{ij}^{(1)}`$ are defined as $`I_{ij}^{(0)}`$ in (44-46) but through $`\delta _1^{(1)}(r)`$, $`\delta _2^{(1)}(r)`$ and $`ϵ^{(1)}(r)`$ instead of $`\delta _1^{(0)}(r)`$, $`\delta _2^{(0)}(r)`$ and $`ϵ^{(0)}(r)`$.
Evidently we cannot consider (48-50) in a manner like (34). But we assume that
$$I_{ij}^{(1)}=I_{ij}^{(0)}+\underset{i,j=1,2,3}{}o(I_{ij}^{(0)}),$$
(51)
where
$$\underset{i,j=1,2,3}{}o(I_{ij}^{(0)})I_{ij}^{(0)},\text{ for }i,j=1,2,3.$$
(52)
This assumption can be considered as a generalization of (39). It is hard to prove in the general case, but our calculations show that this is true at least in case of $`{}_{}{}^{3}SD_1`$ NN scattering.
Eigenphases $`\widehat{\delta }_i^{(0)}`$, $`i=1,2`$ and mixing parameter $`\widehat{ϵ}^{(0)}`$ are real and they define a unitary $`S^{(0)}`$-matrix
$$S^{(0)}=\left(\begin{array}{cc}\mathrm{cos}^2\widehat{ϵ}^{(0)}e^{2ı\widehat{\delta }_{1}^{}{}_{}{}^{(0)}}+\mathrm{sin}^2\widehat{ϵ}^{(0)}e^{2ı\widehat{\delta }_{2}^{}{}_{}{}^{(0)}}& \mathrm{cos}\widehat{ϵ}^{(0)}\mathrm{sin}\widehat{ϵ}^{(0)}\left(e^{2ı\widehat{\delta }_{1}^{}{}_{}{}^{(0)}}e^{2i\widehat{\delta }_{2}^{}{}_{}{}^{(0)}}\right)\\ \multicolumn{2}{c}{}\\ \mathrm{cos}\widehat{ϵ}^{(0)}\mathrm{sin}\widehat{ϵ}^{(0)}\left(e^{2ı\widehat{\delta }_{1}^{}{}_{}{}^{(0)}}e^{2i\widehat{\delta }_{2}^{}{}_{}{}^{(0)}}\right)& \mathrm{sin}^2\widehat{ϵ}^{(0)}e^{2ı\widehat{\delta }_{1}^{}{}_{}{}^{(0)}}+\mathrm{cos}^2\widehat{ϵ}^{(0)}e^{2ı\widehat{\delta }_{2}^{}{}_{}{}^{(0)}}\end{array}\right)$$
(53)
Eigenphases $`\widehat{\delta }_i^{(1)}`$ and mixing parameter $`\widehat{ϵ}^{(1)}`$ are complex but they define $`S^{(1)}`$-matrix in the regular way
$$S^{(1)}=\left(\begin{array}{cc}\mathrm{cos}^2\widehat{ϵ}^{(1)}e^{2ı\widehat{\delta }_{1}^{}{}_{}{}^{(1)}}+\mathrm{sin}^2\widehat{ϵ}^{(1)}e^{2ı\widehat{\delta }_{2}^{}{}_{}{}^{(1)}}& \mathrm{cos}\widehat{ϵ}^{(1)}\mathrm{sin}\widehat{ϵ}^{(1)}\left(e^{2ı\widehat{\delta }_{1}^{}{}_{}{}^{(1)}}e^{2i\widehat{\delta }_{2}^{}{}_{}{}^{(1)}}\right)\\ \multicolumn{2}{c}{}\\ \mathrm{cos}\widehat{ϵ}^{(1)}\mathrm{sin}\widehat{ϵ}^{(1)}\left(e^{2ı\widehat{\delta }_{1}^{}{}_{}{}^{(1)}}e^{2i\widehat{\delta }_{2}^{}{}_{}{}^{(1)}}\right)& \mathrm{sin}^2\widehat{ϵ}^{(1)}e^{2ı\widehat{\delta }_{1}^{}{}_{}{}^{(1)}}+\mathrm{cos}^2\widehat{ϵ}^{(1)}e^{2ı\widehat{\delta }_{2}^{}{}_{}{}^{(1)}}\end{array}\right).$$
(54)
The complex eigenphases $`\widehat{\delta }_i^{(1)}`$ and mixing parameter $`\widehat{ϵ}^{(1)}`$ are defined by the experimental $`S`$-matrix ($`SS^{(1)}`$). A direct consequence of (48-51) is
$$\mathrm{}\delta _i^{(1)}=\delta _i^{(0)},i=1,2;\mathrm{}ϵ^{(1)}=ϵ^{(0)}.$$
(55)
$$\mathrm{}\delta _i^{(1)}=\underset{j=1}{\overset{3}{}}\alpha _jI_{ij}^{(0)},i=1,2;\mathrm{}ϵ_i^{(1)}=\underset{j=1}{\overset{3}{}}\alpha _jI_{3j}^{(0)}.$$
(56)
We may calculate coefficients $`I_{ij}^{(1)}I_{ij}^{(0)}(i,j=1,2,3)`$ using (44-46). A simpler method is to use the following implication of (48-51)
$$I_{ij}^{(1)}=\mathrm{}\frac{\delta _i^{(1)}}{\alpha _j},i=1,2;I_{3j}^{(1)}=\mathrm{}\frac{ϵ^{(1)}}{\alpha _j}.$$
(57)
Calculated coefficients may be checked by (41-43). Next, we calculate $`\alpha _i(i=1,2,3)`$ from (48-50).
## 4 The optical potentials
We apply the developed method of inversion to analysis $`NN`$, $`\pi ^{}N`$ and $`K^+N`$ data up to energies where relativistic effects are essential. We take into account these effects in the frames of relativistic quantum mechanics of systems with a fixed number of particles. The review of this approach can be found in . Here we give only some extracts of it.
The relativistic quantum mechanics of systems consisting of a fixed number of particles is based on the conjecture that the number of particles is constant at not very high energies and on the assumption that the group of invariance for the system under consideration is the Poincare group rather than the Galilei one. A system of two particles is described by the wave function, which is an eigenfunction of the mass operator. In this case we may represent this wave function as a product of the external and internal wave functions . The internal wave function $`\chi `$ is also an eigenfunction of the mass operator and satisfies the following equation
$$\left[\sqrt{\widehat{q}^2+m_1^2}+\sqrt{\widehat{q}^2+m_2^2}+V_{int}\right]\chi =M\chi ,$$
(58)
where $`V_{int}`$ is some interaction operator acting only through internal variables (spins and relative momentum), $`\widehat{q}`$ is a momentum operator of one of the particles in the center of masses frame. Rearrangement of (58) gives
$$\left[\widehat{q}^2+2mV\right]\chi =q^2\chi ,$$
(59)
where
$$q^2=\frac{M^2}{4}\frac{m_1^2+m_2^2}{2}+\frac{(m_1^2m_2^2)^2}{4M^2},$$
(60)
$`m`$ is taken as a nonrelativistic reduced mass
$$m=\frac{m_1m_2}{m_1+m_2},$$
(61)
$`V`$ is an operator acting like $`V_{int}`$ only through internal variables.
In case of two particles with equal masses $`m_1=m_22m`$
$$q^2=\frac{M^2}{4}4m^2.$$
(62)
Eq. (59) is identical in form to the Schrôdinger equation. The quasicoordinate representation corresponds to the realization $`𝐪=i\frac{}{𝐫}`$, $`V=V(𝐫)`$. In we showed that this formalism can be easily generalized for the case of inelastic channels, particularly it allows to take into account isobar channels in NN scattering. This formal coincidence allows us to apply our inversion algorithm.
We applied the described algorithm of inversion to reconstruction of the nucleon-nucleon potential. As input data for this reconstruction we used modern phase shift analysis data up to 1100 MeV for $`{}_{}{}^{3}SD_1`$ state and up to 3 GeV for $`{}_{}{}^{1}S_{0}^{}`$ state of nucleon-nucleon system . The deuteron properties were taken from . These data allow to construct nucleon-nucleon partial potentials sustaining forbidden bound states (Moscow potential introduced in ). Parameters of forbidden bound states for the partial Moscow potentials were chosen to be equal to those of model potentials (see Sect. 2). In this way we constructed the NN optical potentials for $`{}_{}{}^{1}S_{0}^{}`$ and $`{}_{}{}^{3}SD_1`$ partial waves. These potentials describe the deuteron properties and the phase shift analysis data. The $`{}_{}{}^{1}S_{0}^{}`$ phase shifts of Moscow potential begin from $`\pi `$. $`{}_{}{}^{3}S_{0}^{}`$ phase shifts of Moscow potential begin from $`2\pi `$. The mixing parameter $`ϵ_1`$ of Moscow potential differs from that of traditional repulsive core potential by sign. The real parts of the constructed partial potentials are presented in fig. 1. The calculated values of deuteron properties are compared with the experimental data in Table 1. Only three parameters are fixed as input data of inversion problem. These parameters are energy, $`A_S`$ and $`\eta _{d/s}`$. Figs. 1 and 3 demonstrate how changes of $`\delta `$ influence the partial potential for $`{}_{}{}^{1}S_{0}^{}`$ wave.
As another example of application we analyzed the modern $`P31`$ $`\pi ^{}N`$ data up to 2 GeV and $`S01`$ $`K^+N`$ data up to 1 GeV and constructed the corresponding optical potentials. The real parts of the constructed partial potentials for $`P31`$ $`\pi ^{}N`$ data up to 2 GeV and $`S01`$ $`K^+N`$ are presented in fig. 2.
From eqs. (38) and (56, 57) we calculated parameters $`\alpha `$ and $`\alpha _i,(i=1,2,3)`$ which define the imaginary parts of potentials. Our predictions were justified. Calculations with optical and real potentials (fig. 3,7,8) show the validity of (37) and (55). The $`\alpha `$’s predicted by (38) may be improved by a simple numerical method. Predicted and improved values of $`\alpha `$’s are shown in fig. 4,5,6. In all figures ”Calc. I” means calculations from predicted values of (37), ”Calc. II” means calculations from refined values. Fig. 5 shows that calculation of $`\alpha _i,(i=1,2,3)`$ from (56 and 57) does not require refinement because it uses more precise values of $`I_{ij}^{(1)}`$ than those implied by assumption that $`I_{ij}^{(1)}=I_{ij}^{(0)}`$.
In eq. (36) we use parametrization of the $`S`$-matrix from whereas parametrization of partial wave analysis is based on type-$`K`$ scheme . For uncoupled waves, the $`S`$-matrix is given by
$$\begin{array}{c}S=\frac{1K_i+ıK_r}{1+K_iıK_r},\\ \\ \text{where }K_r=tan\stackrel{~}{\delta },K_i=tan^2\stackrel{~}{\rho }.\end{array}$$
(63)
So we had to recalculate data of into $`S`$-matrix, then into parameters of to get input data of inverse problem. Our results are presented in parametrization of .
All potentials and inelasticity multipliers ($`\alpha `$’s) may be downloaded from cite www.physics.khstu.ru in numerical form.
## 5 Conclusions
Let us summarize the results presented in this work. In the first place we mention a presented analytical solution of the Marchenko equation for coupled partial waves in case of diagonal Padé approximant of the corresponding $`S`$-matrix. The inverse scattering scheme at fixed angular momentum is used to construct a local real energy independent potential as a first step of our inversion procedure for single and coupled waves. Furthermore, we consider what changes of $`S`$-matrix are induced by certain transformation of real potential. We have found out that certain simple transformation may have a negligible effect on phase shift but introduce a controllable inelasticity. This transformation does not change the real part of the potential but adds an imaginary part. As a result we get an optical potential with energy independent real part and energy dependent imaginary part. We apply this scheme to NN, $`\pi ^{}N`$ and $`K^+N`$ scattering successfully.
## 6 Appendix
We present only nontrivial integrals of (28). They can be derived from the recursion relations for the Riccati-Hankel functions and from known integrals .
$$\begin{array}{cc}I_1(x,\beta _i,\beta _k,l)=\underset{x}{\overset{\mathrm{}}{}}h_l^+(\beta _it)h_l^+(\beta _k)𝑑t=& \{\begin{array}{c}x(\left(h_l^+\left(\beta _ix\right)\right)^2h_{l1}^+\left(\beta _ix\right)h_{l+1}^+\left(\beta _ix\right))/2,fori=k\\ \frac{\beta _ih_{l1}^+\left(\beta _ix\right)h_l^+\left(\beta _kx\right)\beta _kh_l^+\left(\beta _ix\right)h_{l1}^+\left(\beta _kx\right)}{\beta _{ik}},forik.\end{array}\end{array}$$
(64)
$$\begin{array}{cc}\underset{x}{\overset{\mathrm{}}{}}h_l^+(\beta _it)h_l^+(\beta _kt)t𝑑t=& \{\begin{array}{c}\left(x\left(h_l^+\left(\beta _ix\right)\right)^2+I_1(x,\beta _i,\beta _i,l)\right)/(2\beta _i),fori=k\\ \begin{array}{c}\frac{2\beta _i}{\beta _{ik}^2}\left(\beta _kh_l^+\left(\beta _ix\right)h_{l1}^+\left(\beta _kx\right)\beta _ih_{l1}^+\left(\beta _ix\right)h_l^+\left(\beta _kx\right)\right)+\\ +\frac{1}{\beta _{ik}}h_{l1}^+\left(\beta _ix\right)h_l^+\left(\beta _kx\right)\\ \frac{x}{\beta _{ik}}\left(\beta _kh_l^+\left(\beta _ix\right)h_{l1}^+\left(\beta _kx\right)\beta _ih_{l1}^+\left(\beta _ix\right)h_l^+\left(\beta _kx\right)\right),forik.\end{array}\end{array}\end{array}$$
(65)
Where $`\beta _{ik}=\beta _i^2\beta _k^2`$
$$\begin{array}{cc}\underset{x}{\overset{\mathrm{}}{}}h_l^+(\beta _it)h_l^+(\beta _kt)t^2𝑑t=& \{\begin{array}{c}I_2(\beta _ix,l)/\beta _i^3,fori=k\\ \begin{array}{c}8\frac{\beta _i\beta _k}{\beta _{ik}^3}\left(\beta _kh_l^+\left(\beta _ix\right)h_{l1}^+\left(\beta _kx\right)\beta _ih_l^+\left(\beta _kx\right)h_{l1}^+\left(\beta _ix\right)\right)+\\ +2\frac{\beta _i}{\beta _{ik}^2}(h_l^+\left(\beta _ix\right)h_{l1}^+\left(\beta _kx\right)+\beta _kxh_l^+\left(\beta _ix\right)h_{l1}^+\left(\beta _kx\right)\\ \beta _ixh_l^+\left(\beta _kx\right)h_{l1}^+\left(\beta _ix\right))+\\ +\frac{x}{\beta _{ik}}(h_l^+\left(\beta _kx\right)h_{l1}^+\left(\beta _ix\right)h_l^+\left(\beta _ix\right)h_{l1}^+\left(\beta _kx\right)\\ \beta _kxh_l^+\left(\beta _ix\right)h_{l1}^+\left(\beta _kx\right)+\beta _ixh_l^+\left(\beta _kx\right)h_{l1}^+\left(\beta _ix\right))\\ +\frac{2\beta _k}{\beta _{ik}^2}(h_l^+\left(\beta _kx\right)h_{l1}^+\left(\beta _ix\right)\beta _kxh_l^+\left(\beta _ix\right)h_{l1}^+\left(\beta _kx\right)+\\ +\beta _ixh_l^+\left(\beta _kx\right)h_{l1}^+\left(\beta _ix\right)),forik.\end{array}\end{array}\end{array}$$
(66)
$$I_2(z,1)=\left(ız^2/23z/29ı/2+1/z\right)\mathrm{exp}(ı2z)$$
(67)
$$I_2(z,2)=\left(ız^2/2+7z/2+49ı/224/z24ı/x^2+12/x^3\right)\mathrm{exp}(ı2z)$$
(68)
$$I_2(z,3)=\left(ız^2/213z/2169ı/2+171/z+450ı/x^2765/x^3810ı/x^4+405/x^5\right)\mathrm{exp}(ı2z)$$
(69)
$$\begin{array}{c}I_2(z,4)=(ız^2/2+21z/2+441ı/2745/z3510ı/x^2+11835/x^3+\\ +28560ı/x^447880/x^550400ı/x^6+25200/x^7)\mathrm{exp}(ı2z)\end{array}$$
(70) |
warning/0506/quant-ph0506152.html | ar5iv | text | # Range criterion and classification of true entanglement in 2×𝑀×𝑁 system
## Abstract
We propose a range criterion which is a sufficient and necessary condition satisfied by two pure states transformable with each other under reversible stochastic local operations assisted with classical communication. We also provide a systematic method for seeking all kinds of true entangled states in the $`2\times M\times N`$ system, and can effectively distinguish them by means of the range criterion. The efficiency of the criterion and the method is exhibited by the classification of true entanglement in some types of the tripartite systems.
One of the main tasks in quantum information theory (QIT) is to find out how many different ways there exist, in which several spatially distributed objects could be entangled under certain prior constraints of physical resource. Such a restriction is the transformation of entangled states by local operations assisted with classical communication (LOCC). All bipartite pure entangled states are interconvertible in the asymptotic LOCC transformations Bennett1 . This implies that all bipartite pure entangled states can be used to perform the same task of entanglement processing in the asymptotic regime. However, for single copies of states, any two pure states in the same class under LOCC can be convertible with certainty from each other by local unitary operation Vidal ; Bennett2 .
In the general bipartite system, infinitely many entangled states are not related by local unitary operations and continuous parameters are used to characterize all equivalent classes corresponding to them Linden ; Kempe . In order to classify them more succinctly, an alternative restriction was introduced Bennett2 ; Dur . This restriction requires that the conversion of the states is through stochastic local operations and classical communication (SLOCC), i.e. through LOCC but without imposing that it has to be achieved with certainty. Two pure states of a multipartite system are equivalent under SLOCC if and only if (iff) they are related by an invertible local operator (ILO) Dur . Comparing with the other constraints, SLOCC is coarser, but more practical and simpler for the classification of entanglement. For instance, the state $`|00`$ and the Bell state, $`|\mathrm{\Phi }=|00+|11`$ are under this criterion the only two classes in two-qubit systems.
Such restriction becomes stricter in multipartite settings, since multiparty entanglement has a much more complicated configuration than the bipartite case. For the true tripartite entangled states, there exist not only the three-qubit GHZ state Greenberger and but also the so-called three-qubit W state Dur , which was shown to be essentially different from the GHZ state under SLOCC. In principle, the conclusion of the SLOCC equivalence under ILO’s Dur is sufficient to classify the entanglement properties of single copies of states. However, for the classification of entangled states of general multipartite system, the classification based on this conclusion becomes more and more complicated with increasing of dimensions of Hilbert space of such system and the exponential increase of number of parameters when classifying multiparty entanglement Verstraete1 ; Miyake2 ; Miyake3 . In fact, it is almost impossible to make use of it to perform the classification and the construction of the different entangled states in general multipartite system. So we need to establish a more effective criterion involved in the subspaces of the Hilbert space.
In this paper, we introduce a more straightforward and effective criterion of the equivalent classes of entanglements under SLOCC in general multipartite system, where only a single copy of state is available. Different from the definition of equivalence classes of entangled states suggested by Dür et al Dur , our criterion is constructive for the entangled states of multipartite system. Indeed, based on the criterion, an iterated method can be introduced to determine all classes of true entangled states in the $`2\times M\times N`$ system. Combined such method with the criterion, we not only can classify the entangled states but also can construct them.
In order to arrive at our criterion in a smooth way, we ought to introduce some necessary notations and useful concepts in this paper. Since we are concerned with the entanglement properties of the states of many parties, we should understand the quantum properties of the states of each party. Such quantum state can be described by the reduced density matrix of single party from the composite state of many parties. The rank of it is called as the local rank. Since we are interested with the classification of pure entangled states in the multipartite system, the state $`\rho _{A_1A_2\mathrm{}A_N}`$ can be realized by the expression $`\rho _{A_1A_2\mathrm{}A_N}=|\mathrm{\Psi }_{A_1A_2\mathrm{}A_N}\mathrm{\Psi }|`$. In the notation, we shall use $`\rho _{\mathrm{\Psi }_{AB}}^B`$ to stand for the reduced density operator of $`B`$ system from the state $`|\mathrm{\Psi }_{AB}`$. The local ranks of any pure state are invariant under SLOCC Dur . Hence, one can use the local ranks of the parties to characterize the Hilbert space of multipartite system, e.g., the $`D_1\times D_2\times \mathrm{}\times D_N`$ space, where $`D_i`$ is the local rank of the party $`A_i`$. Due to the invariance of the local ranks under SLOCC and the Schmidt decomposition with respect to the party $`A_j`$ and the other party $`A_1A_2\mathrm{}A_{j1}A_{j+1}\mathrm{}A_N`$, any state $`|\mathrm{\Psi }_{A_1A_2\mathrm{}A_N}`$ ( sometimes also written $`|\mathrm{\Psi }_{D_1D_2\mathrm{}D_N}`$ ) in the $`D_1\times D_2\times \mathrm{}\times D_N`$ space can always be transformed into the following form,
$$|\mathrm{\Phi }=\underset{i=0}{\overset{D_j1}{}}|i_{A_j}|i_{A_1A_2\mathrm{}A_{j1}A_{j+1}\mathrm{}A_N},$$
(1)
where $`i|k_{A_j}=\delta _{ik}`$ and $`\{|i_{A_1A_2\mathrm{}A_{j1}A_{j+1}\mathrm{}A_N},i=0,1,\mathrm{},D_j1\}`$ are a set of linearly independent vectors. We call the above expression $`adjoint`$ $`form`$ and the vector $`|i_{A_1A_2\mathrm{}A_{j1}A_{j+1}\mathrm{}A_N}`$ $`adjoint`$ $`state`$. Here the set $`\{|i_{A_j}\}`$ are chosen as the computational basis.
The concept of range of quantum state plays an essential role in our criterion. Let state $`\rho `$ act on the Hilbert space $``$. In the standard manner, the range of $`\rho `$ is defined by $`R(\rho )=|\mathrm{\Psi }:\rho |\mathrm{\Phi }=|\mathrm{\Psi }`$, for some $`|\mathrm{\Phi }`$. For the $`adjoint`$ reduced density matrix $`\rho _{\mathrm{\Psi }_{A_1A_2\mathrm{}A_N}}^{A_1A_2\mathrm{}A_{j1}A_{j+1}\mathrm{}A_N}`$ of the party $`A_j`$, one can clearly see that all states $`|\mathrm{\Theta }`$, which are from $`|\mathrm{\Theta }=\rho _{\mathrm{\Psi }_{A_1A_2\mathrm{}A_N}}^{A_1A_2\mathrm{}A_{j1}A_{j+1}\mathrm{}A_N}|\mathrm{\Gamma }`$ for any $`|\mathrm{\Gamma }_{A_1A_2\mathrm{}A_{j1}A_{j+1}\mathrm{}A_N}`$, span the whole range of it. In fact, for a general multipartite system, the local rank of each party and the range of the adjoint reduced density matrix of it determine completely the character property of a pure multiple entangled state under SLOCC. The following theorem exhibits such relation.
Theorem 1: Range Criterion. Two pure states of a multipartite system are equivalent under SLOCC iff (i) they have the same local rank of each party, and (ii) the ranges of the adjoint reduced density matrices of each party of them are related by certain ILO’s.
Now, let us use $`|\mathrm{\Psi }_{A_1A_2\mathrm{}A_N}`$ and $`|\mathrm{\Phi }_{A_1A_2\mathrm{}A_N}`$ to denote such two states, and $`V_i`$ to do the ILO acting on the party $`A_i`$. The theorem 1 can be formulated that $`|\mathrm{\Psi }_{A_1A_2\mathrm{}A_N}=V_1V_2\mathrm{}V_N|\mathrm{\Phi }_{A_1A_2\mathrm{}A_N}`$ iff they have the same local rank, $`S_1S_2`$ and $`S_2S_1`$, i.e. $`S_1=S_2`$, where
$`S_1`$ $``$ $`\{|\varphi _{A_2\mathrm{}A_N},|\varphi R(\rho _{\mathrm{\Psi }_{A_1A_2\mathrm{}A_N}}^{A_2\mathrm{}A_N})\},`$ (2)
$`S_2`$ $``$ $`\{|\varphi _{A_2\mathrm{}A_N},|\varphi V_2\mathrm{}V_NR(\rho _{\mathrm{\Phi }_{A_1A_2\mathrm{}A_N}}^{A_2\mathrm{}A_N})\}.`$
Proof. Necessity. If the two states are equivalent under SLOCC, they must have the same local rank. By means of the expression of $`|\mathrm{\Phi }_{A_1A_2\mathrm{}A_N}`$ in the adjoint form, we get $`R(\rho _{\mathrm{\Psi }_{A_1A_2\mathrm{}A_N}}^{A_2\mathrm{}A_N})=_{i,j=0}^{D_11}i\left|V_1^{}V_1\right|j_{A_1}i\left|V_2^{}\mathrm{}V_N^{}\right|\mu _{A_2\mathrm{}A_N}V_2\mathrm{}V_N|j_{A_2\mathrm{}A_N}_{j=0}^{D_11}C_j(V_2\mathrm{}V_N|j_{A_2\mathrm{}A_N})`$, and $`R(\rho _{\mathrm{\Phi }_{A_1A_2\mathrm{}A_N}}^{A_2\mathrm{}A_N})=_{j=0}^{D_11}j|\omega _{A_2\mathrm{}A_N}|j_{A_2\mathrm{}A_N}.`$ Here, $`|\mu `$ and $`|\omega `$ are two arbitrary vectors in the $`A_2\mathrm{}A_N`$ space, and we can write $`|\omega =(\omega _0,\omega _1,\mathrm{},\omega _{d1})^T,d=D_2\times D_3\times \mathrm{}\times D_N`$. We prove that there always exists a vector $`|\omega `$ which satisfies $`C_j=j|\omega ,j=0,1,\mathrm{},D_11`$, namely, $`M_{D_1\times d}(\omega _0,\omega _1,\mathrm{},\omega _{d1})^T=(C_0,C_1,\mathrm{},C_{D_11})^T`$, where the coefficient matrix $`M`$ consists of $`j`$’s index entering the row. Notice that $`\{|i_{A_2\mathrm{}A_N},i=0,1,\mathrm{},D_11\}`$ are linearly independent, so $`D_1d`$. This implies the existence of $`|\omega `$. Since $`V_i`$’s are invertible, the above discussion shows $`S_1S_2`$. Similarly we can get $`S_2S_1`$($`\mathrm{\Phi }\mathrm{\Psi }`$, $`V_iV_i^1`$).
Sufficiency. Suppose $`|\mathrm{\Psi }_{A_1A_2\mathrm{}A_N}`$ and $`|\mathrm{\Phi }_{A_1A_2\mathrm{}A_N}`$ have the same local rank. The above result for $`S_2S_1`$ can be more explicitly expressed as: if we write $`|\mathrm{\Phi }=_{i=0}^{D_11}|i_{A_1}|\phi _i_{A_2\mathrm{}A_N}`$ and $`|\mathrm{\Psi }=_{i=0}^{D_11}|i_{A_1}|\psi _i_{A_2\mathrm{}A_N}`$ , then $`V_2\mathrm{}V_N|\phi _j=_{i=0}^{D_11}a_{ji}|\psi _i,j=0,1,\mathrm{},D_11.`$ This expression can be equivalently written as $`V^{d\times d}|j=_{i=0}^{D_11}a_{ji}|i,j=0,1,\mathrm{},D_11,`$ where the $`|i`$’s are computational basis and $`V^{d\times d}`$ is invertible. So the matrix $`M^{D_1\times D_1}`$ is invertible where $`M_{ij}=a_{ji}`$. Thus, we have $`_{j=0}^{D_11}|j_{A_1}V_2\mathrm{}V_N|\phi _j_{A_2\mathrm{}A_N}=_{j=0}^{D_11}M^T|j_{A_1}|\psi _j_{A_2\mathrm{}A_N}`$. By setting $`V_1=(M^T)^1`$, we arrive at the conclusion of the theorem. For the case of $`S_1S_2`$, the similar result can be proved. Q.E.D.
The above theorem gives a universal criterion of the equivalent classes of multiparty entanglement under SLOCC. It shows that one can judge whether two N-partite states are equivalent under SLOCC or not by analyzing the range of the reduced density operator of any $`N1`$ parties (the ranks are easily obtained). When two given states $`|\mathrm{\Psi }_{A_1A_2\mathrm{}A_N}`$ and $`|\mathrm{\Phi }_{A_1A_2\mathrm{}A_N}`$ are investigated, we have to write out those local operations in detail, i.e., $`V_i^{D_i\times D_i}=[a_{ijk}]`$, the indices of parties, $`i=0,1,\mathrm{},N1,`$ that of the matrix entries $`j,k=0,1,\mathrm{},D_i1`$, and they map one range to another. Obviously, whether the resulting entries $`a_{ijk}`$’s keep the nonsingularity of all $`V_i`$’s determines whether the two states are equivalent under SLOCC or not. To simplify the calculation, one can determine the number of product states in the range of the adjoint reduced density matrix of each party. Such the number of product states can be accounted from the state expansions in the range basic vectors to satisfy the condition of product state. For example, consider the two stand forms of three-qubit states, $`|GHZ_{ABC}=|000+|111`$ and $`|W_{ABC}=|001+|010+|100`$. One can readily write out $`R(\rho _{GHZ}^{AB})=\alpha _0|00+\alpha _1|11`$ , and $`R(\rho _W^{AB})=\beta _0|00+\beta _1(|01+|10)`$, $`\alpha _0,\alpha _1,\beta _0,\beta _1C.`$ Hence, no matter how the coefficients change, there are two product states $`|00`$ and $`|11`$ in $`R(\rho _{GHZ}^{AB})`$, and only one product state $`|00`$ in $`R(\rho _W^{AB})`$, which has been used for the existence of different types of entanglement in three-qubit states Dur . The reason is, the states consisting of adjoint states are always different notation1 . Since any product state can only be transformed into another product state under ILO’s, we obtain a useful corollary as follows.
Corollary 1. If two pure states of a multipartite system are equivalent under SLOCC, the numbers of product states in the ranges of the adjoint reduced density matrices of each party of them must be equal.
The above corollary gives a necessary condition of equivalent classes of multiparty entanglement under SLOCC, and it is by employing this condition often more practical than the analysis of the whole range of given states, e.g., two states $`|GHZ`$ and $`|W`$ are inequivalent under SLOCC, since we have known they have different numbers of product states in the ranges. However, the equality of the numbers of product states in the ranges of two states does not imply the equivalence of them under SLOCC. Once that the equality of the numbers of product states in the ranges is satisfied, one must check whether the states in the ranges can be transformed with each other under ILO’s or not. If yes, such two states belong to an equivalent class under SLOCC. Otherwise, they are not equivalent under SLOCC. We introduce a notation $`[a_1,a_2,a_3,\mathrm{},a_N]`$ representing a set of states, and a state $`|\mathrm{\Phi }_{A_1A_2\mathrm{}A_N}[a_1,a_2,a_3,\mathrm{},a_N]`$ iff the number of product states in $`R(\rho _{\mathrm{\Phi }_{A_1A_2\mathrm{}A_N}}^{A_1\mathrm{}A_{i1}A_{i+1}\mathrm{}A_N})`$ is $`a_i`$, $`i=2,\mathrm{},N1`$. We will characterize the entanglement property of the state $`|\mathrm{\Phi }_{A_1A_2\mathrm{}A_N}`$ mainly by this notation.
We examine the above depiction with a concrete example. Consider two states in the $`2\times 4\times 4`$ system, $`|\varphi _0_{ABC}=|000+|111+|022+|033,`$ and $`|\varphi _1_{ABC}=|001+|010+|100+|022+|033.`$ Notice both of them possess symmetry under exchange of particle $`B`$ and $`C`$, so we can write $`R(\rho _{\varphi _0}^{AB})=R(\rho _{\varphi _0}^{AC})=\alpha _0|00+\alpha _1|11+\alpha _2|02+\alpha _3|03`$, and $`R(\rho _{\varphi _0}^{BC})=\beta _0(|00+|22+|33)+\beta _1|11,\alpha _0,\alpha _1,\alpha _2,\alpha _3,\beta _0,\beta _1C`$. Let us move to find out the number of product states in all three ranges. one can readily see there are infinitely many product states in $`R(\rho _{\varphi _0}^{AB})`$ (also $`R(\rho _{\varphi _0}^{AC})`$), e.g., by supposing $`\alpha _0=\alpha _1=0`$ and then changing $`\alpha _2`$ and $`\alpha _3`$. On the other hand, supposing $`\beta _0=0`$ shows $`|11`$ is the unique product state in $`R(\rho _{\varphi _0}^{BC})`$. According to the above notation, we then write
$$|\varphi _0_{ABC}=|000+|111+|022+|033[1,\mathrm{},\mathrm{}].$$
In a similar way, we obtain
$$|\varphi _1_{ABC}=|001+|010+|100+|022+|033[1,\mathrm{},\mathrm{}].$$
Clearly, the numbers of product states in each range of $`|\varphi _0`$ and $`|\varphi _1`$ are equal. For judging whether they are equivalent under SLOCC or not, we have to analyze the detailed structure of range of them. First, write out $`R(\rho _{\varphi _1}^{BC})=\beta _0(|01+|10+|22+|33)+\beta _1|00`$, whose local rank of system B (or C) is 4 or 1. According to the range criterion, it must be mapped into $`R(\rho _{\varphi _0}^{BC})`$ by some ILO’s. Because the local rank of system B (or C) in $`R(\rho _{\varphi _0}^{BC})`$ is 3 or 1, there exist no ILO’s making this transformation. $`|\varphi _0`$ and $`|\varphi _1`$ are thus inequivalent under SLOCC, although they have the same number of product states in each range.
Now, let us focus on the classification of entanglement in the $`2\times M\times N`$ system. Although the structure of the whole space of pure states of this system $`|\mathrm{\Psi }_{2\times M\times N}`$ is very complicated, we can establish an interesting relation between $`|\mathrm{\Psi }_{2\times M\times N}`$ and some spaces of pure states with the local ranks lower than those of it. We shall use $``$ to denote the equivalence under SLOCC in this paper. There are two useful ILO’s in the below discussion. They are defined by $`O_1^A(|\varphi ,\alpha ):|\varphi _A\alpha |\varphi _A`$ and $`O_2^A(|\varphi ,|\psi ):|\varphi _A|\varphi _A+|\psi _A`$, respectively.
Any state $`|\mathrm{\Phi }_{ABC}`$, in the $`2\times M\times N`$ space for $`2MN2M`$, can be expanded as its adjoint form $`|\mathrm{\Phi }_{ABC}=_{i=0}^{N1}|\psi _i_{AB}|i_C`$, where $`\{|\psi _i_{AB},i=0,\mathrm{},N1\}`$ are linearly independent and $`\{|i_C,i=0,\mathrm{},N1\}`$ are a set of computational basis. According to the conclusion of the lemma 10 in Kraus , we know that there is always at least one product state in $`R(\rho _{\mathrm{\Phi }_{ABC}}^{AB})`$. Suppose $`_{i=0}^{N1}c_i|\psi _i_{AB}`$ is of product form, where the constants $`c_i`$’s do not equal zero simultaneously. Let $`c_k0,k[0,N1]`$. By performing the operations $`O_1^C(|k,c_k)_{i=0,ik}^{N1}O_2^C(|i,c_i|k)`$, $`|k_C|N1_C`$ and some ILO’s making $`_{i=0}^{N1}c_i|\psi _i_{AB}|0,M1_{AB}`$, one can obtain $`|\mathrm{\Phi }_{ABC}_{i=0,ik}^{N1}|\psi _i_{AB}|i_C+_{i=0}^{N1}c_i|\psi _i_{AB}|k_C|0,M1,N1+_{i=0}^{N2}|\psi _i^{}_{AB}|i_C`$. Thus $`\{|\psi _i^{}_{AB},i=0,\mathrm{},N2\}`$ remains a set of linearly independent vectors. Let $`|\mathrm{\Psi }=_{i=0}^{N2}|\psi _i^{}_{AB}|i_C`$, we have
$`M`$ $`=`$ $`rank(\rho _{\mathrm{\Phi }_{ABC}}^B)=rank(|M1M1|+\rho _{\mathrm{\Psi }_{ABC}}^B)`$
$`rank(|M1M1|)+rank(\rho _{\mathrm{\Psi }_{ABC}}^B),`$
so $`rank(\rho _{\mathrm{\Psi }_{ABC}}^B)M1`$. Based on the restriction of local ranks, we obtain a general equivalence relation such that $`|\mathrm{\Phi }_{ABC}|0,M1,N1+|\mathrm{\Psi }^{}_{ABC}`$, $`|\mathrm{\Psi }^{}_{ABC}`$ lies in $`(|\mathrm{\Psi }_0_{2\times (M1)\times (N1)},|\mathrm{\Psi }_1_{2\times M\times (N1)})`$ and $`(|\mathrm{\Psi }_2_{1\times (M1)\times (N1)},|\mathrm{\Psi }_3_{1\times M\times (N1)})`$. That is, any state $`|\mathrm{\Psi }_{2\times M\times N}`$ can be equivalently transformed into one of the four kinds of states by some ILO’s. We are going to simplify this relation so that it is more efficient. Let $`2MN`$, without loss of generality. If $`M=N`$, the ILO’s make $`|0,M1,N1+|\mathrm{\Psi }_2|0,M1,N1+|\mathrm{\Psi }_3|0,M1,M1+|1_{i=0}^{M2}|ii|1,M1,M1+(|000+|1_{i=1}^{M2}|ii)|0,M1,N1+|\mathrm{\Psi }_0`$. So we only calculate $`|0,M1,N1+|\mathrm{\Psi }_0`$ and $`|0,M1,N1+|\mathrm{\Psi }_1`$. If $`M<N`$, since $`|0,M1,N1+|\mathrm{\Psi }_2`$ leads to a class with lower rank and so does $`|0,M1,N1+|\mathrm{\Psi }_3`$ except the case $`N1=M`$, we do not consider them for true entangled states. The situation of $`N1=M`$ is fully similar to that of $`N=M`$ as above. Now, we arrive at the following conclusion.
Lemma 1. For the classification of true tripartite entangled states, there exists a general equivalence relation under SLOCC such that
$`|\mathrm{\Psi }_{2\times M\times N}|0,M1,N1+\{\begin{array}{c}|\mathrm{\Psi }_0_{2\times (M1)\times (N1)};\hfill \\ |\mathrm{\Psi }_1_{2\times M\times (N1)}.\hfill \end{array}`$
By extracting further the adjoint product state of the $`B`$ party from $`|\mathrm{\Psi }_1`$ and using of the above lemma, we obtain
$`|0,M1,N1+|\mathrm{\Psi }_1_{2\times M\times (N1)}|0,M1,N1`$
$`+|1,M1|\chi +\{\begin{array}{c}|\mathrm{\Phi }_0_{2\times (M1)\times (N1)};\hfill \\ |\mathrm{\Phi }_1_{2\times (M1)\times (N2)},\hfill \end{array}`$
where $`|\chi _{i=0}^{N2}a_i|i`$ and the arbitrary constants $`a_i`$’s do not equal zero simultaneously. By combining the above result with the lemma, we can write out an united relation of equivalence.
Corollary 2. For the classification of true tripartite entangled states under SLOCC, the following equivalence relation is true,
$`|\mathrm{\Psi }_{2\times M\times N}\{\begin{array}{c}|\mathrm{\Omega }_0(a|0+b|1)|M1,N1\hfill \\ +|\mathrm{\Psi }_{2\times (M1)\times (N1)},\hfill \\ |\mathrm{\Omega }_1|0,M1,N1\hfill \\ +|1,M1,N2+|\mathrm{\Psi }_{2\times (M1)\times (N2)},\hfill \\ |\mathrm{\Omega }_2|\mathrm{\Omega }_0+|0,M1|\chi ,b0,\hfill \\ |\mathrm{\Omega }_3|\mathrm{\Omega }_0+|1,M1|\chi ,a0.\hfill \end{array}`$
The condition $`a0`$ or $`b0`$ keeps $`|\mathrm{\Omega }_2`$ and $`|\mathrm{\Omega }_3`$ not becoming $`|\mathrm{\Omega }_0`$.
Such equivalence relation shows that the lower rank classes of the entangled states can be used to generate the higher rank classes of the true entangled states for any $`2\times M\times N`$ system, called as ”Low-to-High Rank Generating Mode” or LHRGM for short. So the corollary and the range criterion of the theorem 1 provide a systematic method to classify all kinds of true tripartite entangled states in the $`2\times M\times N`$ system.
First of all, we re-derive the result of true three qubit entanglements in the present formulation. For the $`2\times 2\times 2`$ system, from corollary 2, $`|\mathrm{\Omega }_0_{2\times 2\times 2}(a|0+b|1)|11+|000`$. It is transformed into $`|111+|000`$ by the ILO $`O_2^A(|1,ab^1|0)`$ for $`b0`$. Hence, $`|\mathrm{\Omega }_0_{2\times 2\times 2}|GHZ`$. By means of $`|0_A|1_A`$, one can easily see that $`|\mathrm{\Omega }_2|\mathrm{\Omega }_3`$. So we only calculate $`|\mathrm{\Omega }_2_{2\times 2\times 2}(a|0+b|1)|11+(c|0+d|1)|00+|010`$. By performing the same ILO’s on $`|\mathrm{\Omega }_0_{2\times 2\times 2}`$, we can transform $`|\mathrm{\Omega }_2_{2\times 2\times 2}`$ into $`|111+(c^{}|0+d^{}|1)|00+|010`$. If $`c^{}0`$, we find that $`|\mathrm{\Omega }_2_{2\times 2\times 2}|GHZ`$ by the ILO’s $`O_2`$ acting on each party. On the other hand, if $`c^{}=0`$, one readily gets $`|\mathrm{\Omega }_2_{2\times 2\times 2}|W`$ by $`|0_A|1_A`$ and $`|0_B|1_B`$. $`|GHZ`$ and $`|W`$ are characterized as two states belonging to the classes $`[2,2,2]`$ and $`[1,1,1]`$, respectively. According to the range criterion of the theorem 1, they are two inequivalent states of true tripartite entanglement under SLOCC.
Now we turn to a more complicated system, the $`2\times 3\times 3`$ system. It can be easily seen from the corollary that all possible classes of true entanglement in this system can be written as the following forms
$`|022+V_AV_BV_C|\mathrm{\Psi }_{2\times 2\times 2},|022+|120+V_AV_BV_C|\mathrm{\Psi }_{2\times 2\times 2},|022+|121+|000+|110,`$
where $`V_A,V_B,V_C`$ are $`2\times 2`$ ILO’s.
Based on the method of the LHRGM, we can use the lower rank entangled classes $`|GHZ`$ and $`|W`$ to generate some entangled classes of the $`2\times 3\times 3`$ system. By using of the explicit forms of $`2\times 2`$ ILO’s, the above expressions can be equivalently written as
$`(a|0+b|1)|22+|000+|111,(I)(a|0+b|1)|22+|001+|010+|100,(II)(a|0+b|1)|22+(c|0+d|1)|2(f|0+g|1)+|000+|111,(III)(a|0+b|1)|22+(c|0+d|1)|2(f|0+g|1)+|001+|010+|100,(IV)|022+|121+|000+|110.(V)`$
Here, $`(a,b),(c,d),(f,g)`$ are arbitrary constants and two constants in each bracket cannot equal zero simultaneously. By classifying the expressions (I), $`\mathrm{}`$, (V) into all possible entangled classes under ILO’s, one can find out all inequivalent classes of true entanglement in the $`2\times 3\times 3`$ system. The result is following.
Theorem 2. There are six classes of true entangled states under SLOCC in the $`2\times 3\times 3`$ system as following
$`|\mathrm{\Psi }_1=|000+|111+(|0+|1)|22[0,3,3],|\mathrm{\Psi }_2=|010+|001+|112+|121[0,\mathrm{},\mathrm{}],|\mathrm{\Psi }_3=|000+|111+|022[1,\mathrm{},\mathrm{}],|\mathrm{\Psi }_4=|100+|010+|001+|112+|121[0,1,1],|\mathrm{\Psi }_5=|100+|010+|001+|022[1,\mathrm{},\mathrm{}],|\mathrm{\Psi }_6=|100+|010+|001+|122[0,2,2].`$
Proof. According to the range criterion, we only have to judge state $`|\mathrm{\Psi }_3`$ and $`|\mathrm{\Psi }_5`$. We notice the product states are $`|11`$ in $`R(\rho _{\mathrm{\Psi }_3}^{BC})`$ and $`|00`$ in $`R(\rho _{\mathrm{\Psi }_5}^{BC})`$. This implies that if they are transformable, $`|0_5^C`$ must be transformed into $`|1_3^C`$. Because the adjoint state of $`|0_5^C`$ is $`|01+|10`$, and the adjoint state of $`|1_3^C`$ is $`|11`$, so there is no ILO between these two states.
Subsequently, we have to prove that the expressions (I), $`\mathrm{}`$, (V) can be only and just transformed into the states in the theorem 2 by the ILO’s. First we observe the expression (I). If $`ab=0`$, by $`|0|1`$ in all systems and $`O_1^B(|2,\alpha )`$ ( $`\alpha `$ is regarded as the possible $`a^1`$, $`b^1`$, etc ), it leads to the state $`|\mathrm{\Psi }_3`$. On the other hand, if $`ab0`$, the operation $`O_1^A(|1,\alpha ^1)O_1^B(|1,\alpha )`$ makes the expression (I) into the state $`|\mathrm{\Psi }_1`$. Second, the expression (II), if $`b=0`$, becomes $`|\mathrm{\Psi }_5`$. For $`b0`$, by means of $`O_2^A(|1,\alpha |0)O_2^C(|1,\alpha |0)`$, we can transform (II) into the state $`|\mathrm{\Psi }_6`$. The ILO’s $`O_2`$ and $`O_1`$ acting on each party produce that the repression (IV) $`|122+|02(f^{}|0+g^{}|1)+|001+|010+|100`$ for $`b0`$ and $`|022+|121+|001+|010+|100`$ for $`b=0`$. Furthermore, by the ILO’s of $`O_2^C`$, $`|0_B|1_B`$, $`|0_C|1_C`$ for $`b0`$ and of $`|1_B|2_B`$, $`|1_C|0_C`$, $`|1_A|0_A`$ for $`b=0`$, we establish the relation the repression (IV) $`|\mathrm{\Psi }_6`$ for $`b0`$ and $`|\mathrm{\Psi }_4`$ for $`b=0`$. The expression (V) leads to state $`|\mathrm{\Psi }_2`$ by $`|0_B|2_B`$ and $`|1|0`$ in all parties. Finally, we have checked that the expression (III) is transformed into the expressions of (I) and (II) under some ILO’s. Q.E.D.
In the $`2\times M\times 2M`$ system, since the dimension of the Hilbert space of the $`AB`$ part is equal to that of the $`C`$ part in the Schmidt decomposition, all pure states of this system are transformed into an unique equivalent class of true entanglement $`|\mathrm{{\rm Y}}_0|0_{i=0}^{M1}|ii+|1_{i=0}^{M1}|i,i+M[0,0,\mathrm{}]`$ under the ILO’s. By using of this result and the method of LHRGM, we can construct the equivalent classes of true entanglement in a general type of the $`2\times M\times N`$ system.
Theorem 3. There are two classes of states under SLOCC in any $`2\times (M+1)\times (2M+1)`$ system ($`M1`$),
$`|\mathrm{{\rm Y}}_1|0,M,2M+|\mathrm{{\rm Y}}_0[0,1,\mathrm{}]`$;
$`|\mathrm{{\rm Y}}_2|0,M,2M+|1,M,M1+|\mathrm{{\rm Y}}_0[0,0,\mathrm{}]`$.
Proof. Following the LHRGM, from corollary 2, we can read out $`|\mathrm{\Omega }_0(a|0+b|1)|M,2M+|\mathrm{{\rm Y}}_0`$. If ab=0, $`|\mathrm{\Omega }_0`$ just is $`|\mathrm{{\rm Y}}_1`$. For the case of $`ab0`$, we perform the ILO $`U_1=O_2^A(|0,\alpha |1)_{i=0}^{M1}O_2^C(|i+M,\alpha |i)`$ on $`|\mathrm{\Omega }_0`$ to be reduced to $`|\mathrm{{\rm Y}}_1`$. In the below proof, the invariance of $`|\mathrm{{\rm Y}}_0`$ under two ILO’s is very useful. They are given by ILO I: $`|0_A|1_A`$, $`|i_C|i+M_C,`$ $`i=0,\mathrm{},M1,`$ and ILO II: $`|i_B|j_B`$, $`|i_C|j_C`$, $`|i+M_C|j+M_C,i,j\{0,\mathrm{},M1\}`$. By the ILO I’s invariance of $`|\mathrm{{\rm Y}}_0`$, it is obvious that $`|\mathrm{\Omega }_2|\mathrm{\Omega }_3(a|0+b|1)|M,2M+|1,M|\chi +|\mathrm{{\rm Y}}_0`$. By using the ILO’s $`U_1`$, $`U_2=_{i=0}^{M1}O_2^B(|i,a_i|M)O_2^C(|2M,_{i=0}^{M1}a_i|i)`$, and the modified $`U_2`$ to act on $`|\mathrm{\Omega }_3`$ sequently, we can transform $`|\mathrm{\Omega }_3`$ into $`|\mathrm{{\rm Y}}_2`$ by means of the ILO II’s invariance of $`|\mathrm{{\rm Y}}_0`$.
Finally we determine the $`|\mathrm{\Omega }_1`$’s family by induction. First of all, let us notice that $`|\mathrm{\Omega }_1_{2\times (M+2)\times (2M+3)}|0,M+1,2M+2+|1,M+1,2M+1+|\mathrm{{\rm Y}}_i_{2\times (M+1)\times (2M+1)}`$ leads to the states $`|\mathrm{{\rm Y}}_i_{2\times (M+2)\times (2M+3)}`$, by both $`|M+1_B|M_B`$ and $`|2M+2_C|2M_C`$. This iterated relation implies that if there are only two classes in $`|\mathrm{\Psi }_{2\times (M+1)\times (2M+1)}`$, $`|\mathrm{\Omega }_1_{2\times (M+2)\times (2M+3)}`$ must enter into two classes of states in $`|\mathrm{\Psi }_{2\times (M+2)\times (2M+3)}`$. For the case $`M=1`$, the conclusion of Miyake2 told us that indeed there exist only two classes of true entangled states in the $`2\times 2\times 3`$ system. Hence, there exist only two inequivalent classes of $`|\mathrm{\Psi }_{2\times (M+1)\times (2M+1)}`$ under SLOCC, and $`|\mathrm{\Omega }_1`$ must be classified into such two classes $`|\mathrm{{\rm Y}}_1`$ and $`|\mathrm{{\rm Y}}_2`$. Q.E.D.
Subsequently, we can use the entangled classes in the theorem 3 to generate those in another type of systems in the LHRGM way. The structure of the entangled classes in this type of system is different from that in theorem 3. We give them here.
Theorem 4. In any $`2\times (M+2)\times (2M+2)`$ system($`M2`$), there are six classes of true entangled states under SLOCC
$`|\mathrm{\Theta }_0|1,M+1,2M+1+|\mathrm{{\rm Y}}_1[0,2,\mathrm{}];|\mathrm{\Theta }_1|0,M+1,2M+1+|\mathrm{{\rm Y}}_1[0,\mathrm{},\mathrm{}];|\mathrm{\Theta }_2|1,M+1,2M+1+|\mathrm{{\rm Y}}_2[0,1,\mathrm{}];|\mathrm{\Theta }_3|0,M+1,2M+1+|1,M+1,2M+|\mathrm{{\rm Y}}_1[0,1,\mathrm{}];|\mathrm{\Theta }_4|0,M+1,2M+1+|1,M+1,0+|\mathrm{{\rm Y}}_2[0,0,\mathrm{}];|\mathrm{\Theta }_5|0,M+1,2M+1+|1,M+1,2M+|\mathrm{{\rm Y}}_2[0,0,\mathrm{}].`$
Proof. According to corollary 2 and the rule of LHRGM, first, we write out the states to be dealt with
$`|\mathrm{\Omega }_0(a|0+b|1)|M+1,2M+1+|\mathrm{{\rm Y}}_j,`$
$`|\mathrm{\Omega }_2|\mathrm{\Omega }_0+|0,M+1_{i=0}^{2M}a_i|i,`$
$`|\mathrm{\Omega }_3|\mathrm{\Omega }_0+|1,M+1_{i=0}^{2M}a_i|i,`$
where $`j=1,2`$. Now, we have to prove that the above states must be transformed into six kinds of states in the theorem 4 by some ILO’s. These ILO’s are composed of those similar to $`U_1`$, $`U_2`$ appeared in the proof of the theorem 3, and the ILO’s of state exchanges in each party. By means of these ILO’s, we have proven the following conclusions. For the case of generation in the direction of the class $`|\mathrm{{\rm Y}}_1`$, it can be easily proved that, by ILO’s $`\stackrel{~}{U_1}=O_2^B(|M,a_{2M}|M+1)`$, and $`|M+1_B|M_B,|2M+1_C|2M_C`$, we have $`|\mathrm{\Omega }_2|\mathrm{\Omega }_0`$. By ILO $`\stackrel{~}{U_2}=O_2^A(|1,\alpha |0)_{i=0}^{M1}O_2^C(|i,\alpha |i+M)`$, $`|\mathrm{\Omega }_0`$ is transferred to $`|\mathrm{\Theta }_0`$ for $`b0`$, and $`|\mathrm{\Theta }_1`$ for $`b=0`$ respectively. Meanwhile, by $`\stackrel{~}{U_3}=O_2^A(|1,\alpha |0)O_2^C(|2M+1,_{i=0}^{2M}a_i|i)O_2^B(|M,\alpha |M+1)`$, and $`|M+1_B|M_B,|2M+1_C|2M_C`$, we obtain $`|\mathrm{\Omega }_3|\mathrm{\Omega }_0`$ for $`b0`$. For $`b=0`$, $`|\mathrm{\Omega }_3|\mathrm{\Omega }_0`$ if there is at least one non-vanishing $`a_i,i=0,1,\mathrm{},M1`$, and $`|\mathrm{\Omega }_3|\mathrm{\Theta }_3`$ if all $`a_i`$’s equal zero, $`i=0,1,\mathrm{},M1`$.
On the other hand, we consider the generation in the direction of the class $`|\mathrm{{\rm Y}}_2`$. It is useful that the invariance of $`|\mathrm{{\rm Y}}_2`$ under the ILO III: $`|0_A|1_A,|M_B|M1_B,|2M_C|2M1_C`$ and $`|i_C|i+M_C,i=0,\mathrm{},M2`$. By virtue of this invariance and the ILO’s, we find that $`|\mathrm{\Omega }_0|\mathrm{\Theta }_2`$. Furthermore, $`|\mathrm{\Omega }_2|\mathrm{\Omega }_3|\mathrm{\Theta }_4`$ if there is at least one non-vanishing $`a_i,i=0,1,\mathrm{},M2`$, and $`|\mathrm{\Omega }_2|\mathrm{\Omega }_3|\mathrm{\Theta }_5`$ if all $`a_i`$’s equal zero, $`i=0,1,\mathrm{},M2`$.
Similarly, one can calculate $`|\mathrm{\Omega }_1`$ by induction. By calculation of $`|\mathrm{\Omega }_0`$, $`|\mathrm{\Omega }_2`$ and $`|\mathrm{\Omega }_3`$ above, the classes in $`|\mathrm{\Psi }_{2\times 3\times 4}`$ are found such that
$`|\mathrm{\Theta }_0_{2\times 3\times 4}|123+|012+|000+|101,|\mathrm{\Theta }_1_{2\times 3\times 4}|023+|012+|000+|101,|\mathrm{\Theta }_2_{2\times 3\times 4}|123+|012+|110+|000+|101,|\mathrm{\Theta }_3_{2\times 3\times 4}|023+|122+|012+|000+|101,|\mathrm{\Theta }_5_{2\times 3\times 4}|023+|122+|012+|110+|000+|101.`$
Notice that the class $`|\mathrm{\Theta }_4`$ disappears here, for the coefficients $`a_i,i=0,\mathrm{},M2`$ always equal zero in the above derivation of $`|\mathrm{\Theta }_4`$ and $`|\mathrm{\Theta }_5`$. Then $`|\mathrm{\Omega }_1_{2\times 3\times 4}|023+|122+|\mathrm{\Psi }_{2\times 2\times 2}`$. For the case of $`|\mathrm{\Psi }_{2\times 2\times 2}|GHZ`$, we have $`|\mathrm{\Omega }_1|023+|122+|000+|111`$. By $`|3_C|1_C`$ and $`|2_B|1_B`$, one obtains $`|\mathrm{\Omega }_1|123+|\mathrm{\Psi }_{2\times 2\times 3}|\mathrm{\Omega }_0`$. On the other hand, if $`|\mathrm{\Psi }_{2\times 2\times 2}|W`$, then $`|\mathrm{\Omega }_1|023+|122+|001+|010+|100`$. By the operations $`|3_C|1_C`$ and $`|2_B|0_B`$, we get $`|\mathrm{\Omega }_1|001+|102+|023+|010+|120|023+|120+|\mathrm{\Psi }_{2\times 2\times 3}|\mathrm{\Omega }_3`$. So there are five classes of entanglement in the $`2\times 3\times 4`$ system. Next, one should write out all classes in $`|\mathrm{\Psi }_{2\times 4\times 6}`$, which is really the first step of the induction. Following the above technique we obtain $`|\mathrm{\Theta }_i_{2\times 4\times 6},i=0,1,2,3,4,5`$. Then we calculate $`|\mathrm{\Omega }_1_{2\times 4\times 6}|035+|134+|\mathrm{\Psi }_{2\times 3\times 4}`$, where $`|\mathrm{\Psi }_{2\times 3\times 4}|\mathrm{\Theta }_i_{2\times 3\times 4},i=0,1,2,3,5.`$ By $`|5_C|3_C`$ and $`|3_B|2_B`$, for the case of $`|\mathrm{\Theta }_i_{2\times 3\times 4},i=0,1,2`$, $`|\mathrm{\Omega }_1|135+|\mathrm{\Psi }_{2\times 3\times 5}|\mathrm{\Omega }_0`$; for the case of $`|\mathrm{\Theta }_i_{2\times 3\times 4},i=3,5`$, $`|\mathrm{\Omega }_1|035+|132+|\mathrm{\Psi }_{2\times 3\times 5}|\mathrm{\Omega }_3`$. So there are six classes of entanglement in the $`2\times 4\times 6`$ system, i.e., $`|\mathrm{\Theta }_i_{2\times 4\times 6},i=0,1,2,3,4,5`$. Similar to the first step, one can continue with the deduction. That is, by $`|2M+3_C|2M+1_C`$ and $`|M+2_B|M+1_B`$, we always get $`|\mathrm{\Omega }_1_{2\times (M+3)\times (2M+4)}|0,M+2,2M+3+|1,M+2,2M+2+|\mathrm{\Psi }_{2\times (M+2)\times (2M+2)}|\mathrm{\Omega }_i,i=0,2,3.`$ Thus, all classes of $`|\mathrm{\Omega }_1`$ belong to $`|\mathrm{\Theta }_i,i=0,1,2,3,4,5.`$
Finally we should check that these classes of states in $`|\mathrm{\Psi }_{2\times (M+2)\times (2M+2)}`$ are inequivalent. According to the range criterion, we need only to discuss the relations between $`|\mathrm{\Theta }_2`$ and $`|\mathrm{\Theta }_3`$, and between $`|\mathrm{\Theta }_4`$ and $`|\mathrm{\Theta }_5`$. The proof of the former case is similar to that in theorem 2. Since the requirement of $`|2M_{\mathrm{\Theta }3}^C|2M+1_{\mathrm{\Theta }2}^C`$ always leads to an entangled adjoint state of $`|2M+1_{\mathrm{\Theta }2}^C`$, $`|\mathrm{\Theta }_2`$ and $`|\mathrm{\Theta }_3`$ are inequivalent. The proof of the latter case is some difficult. We can prove it by the reduction to absurdity. First, we suppose that the AB system is in the adjoint state and there exist the possible ILO’s $`V_A`$ and $`V_B`$ to make $`|\mathrm{\Theta }_4|\mathrm{\Theta }_5`$. By theorem 1, $`V_AV_B[R(\rho _{\mathrm{\Theta }4}^{AB})]R(\rho _{\mathrm{\Theta }5}^{AB})`$. This relation produces a set of equations satisfied by the matrix elements of $`V_A`$ and $`V_B`$. Carefully analyzing this set of equations, we have found that there exist only some singular solutions of $`V_A`$ and $`V_B`$, i.e. Det$`[V_B]=0`$ or Det$`[V_A]=0`$. So there exists no ILO making $`|\mathrm{\Theta }_4`$ and $`|\mathrm{\Theta }_5`$ equivalent. We have detailedly proven this fact in the Appendix. Q.E.D.
To summarize, by using of the range criterion and the method of LHRGM developed by us here, we have finished the construction of true entangled classes of some types of the $`2\times M\times N`$ system, which are with finite kinds of states. For the system with higher dimensions, one can also find out the classification by the techniques in this paper notation2 . However, our method can be also applied to classify the entangled system with infinite kinds of states, which is also a puzzling issue in quantum information theory. Our results are helpful to classify and construct the inequivalent classes of entangled states in many-qubit system.
The work was partly supported by the NNSF of China Grant No.90503009 and 973 Program Grant No.2005CB724508.
APPENDIX: THE INEQUIVALENCE OF $`|\mathrm{\Theta }_4`$ AND $`|\mathrm{\Theta }_5`$
Suppose that the AB system is in the adjoint state and the possible ILO’s are taken as
$$V_A^{2\times 2}=\left(\begin{array}{cc}w& x\\ y& z\end{array}\right),V_B^{(M+2)\times (M+2)}=[a_{ij}],i,j=0,\mathrm{},M+1.$$
Thus, these ILO’s make $`|\mathrm{\Theta }_4|\mathrm{\Theta }_5`$, if det$`[V_A]`$det$`[V_B]0`$ is satisfied. According to theorem 1, we must have
$$V_AV_B[R(\rho _{\mathrm{\Theta }4}^{AB})]R(\rho _{\mathrm{\Theta }5}^{AB}),$$
$`(A1)`$
which yields that ($`M>2`$)
$`V_AV_B[c_0^{}|0,M+|0_{i=1}^{M2}c_i^{}|i+c_{M1}^{}|0,M+1+|1_{i=M}^{2M1}c_i^{}|i+c_{2M}^{}(|1,M+1+|0,0)+c_{2M+1}^{}(|1,M+|0|M1)]_{i=0}^{M2}c_i|0,i+c_{M1}|0,M+1+_{i=M}^{2M1}c_i|1,i+c_{2M}(|1,M+1+|0,M)+c_{2M+1}(|1,M+|0,M1),`$ the coefficients $`c_i^{},c_i`$’s, i=0,…,2M+1 are arbitrarily decided by theorem 1. Let only one $`c_i^{}`$ be nonzero, e.g., $`c_{M1}^{}0`$ and $`c_i^{}=0,iM1`$. The above expression then implies
$$ya_{_{M+1,M+1}}=wa_{_{M,M+1}},ya_{_{M,M+1}}=wa_{_{M1,M+1}}.$$
$`(A2)`$
Notice we have deserted the trivial results that can’t bring those similar to the above relations between the entries of $`V_A`$ and $`V_B`$ without $`c_i`$’s. We continue in the same vein, that is, to choose the uniquely nonzero $`c_i^{}`$ in turn and obtain a sequence of relations such that
$$\{\begin{array}{cc}ya_{_{M+1,i}}=wa_{_{M,i}},i=1,\mathrm{},M2,M,M+1,& \\ ya_{_{M,i}}=wa_{_{M1,i}},i=1,\mathrm{},M2,M,M+1,& \\ za_{_{M+1,i}}=xa_{_{M,i}},i=0,\mathrm{},M1,& \\ za_{_{M,i}}=xa_{_{M1,i}},i=0,\mathrm{},M1,& \\ za_{_{M+1,M}}+ya_{_{M+1,M1}}=xa_{_{M,M}}+wa_{_{M,M1}},& \\ za_{_{M,M}}+ya_{_{M,M1}}=xa_{_{M1,M}}+wa_{_{M1,M1}}.& \end{array}$$
$`(A3)`$
If $`wxyz0`$, it must be that $`a_{_{M1,i}}=a_{_{M,i}}=a_{_{M+1,i}}=0,i=1,\mathrm{},M2,`$ since $`wzxy`$. By substituting the former four expressions of (A3) into the last two expressions of it, on the other hand, one can get $`a_{_{M,M}}/y=a_{_{M,M1}}/z,a_{_{M,M}}/w=a_{_{M,M1}}/x`$. This implies that $`a_{_{M,M}}=a_{_{M,M1}}=0`$, which leads to $`a_{_{M1,i}}=a_{_{M,i}}=a_{_{M+1,i}}=0,i=1,\mathrm{},M.`$ This conclusion is just det$`[V_B]=0`$. So there must be at least one equaling zero in $`\{w,x,y,z\}.`$ If $`x=0,wz0`$ (or $`y=0,wz0`$,etc), it will lead to $`a_{_{M,i}}=a_{_{M+1,i}}=0,i=0,\mathrm{},M,`$ again det$`[V_B]=0`$. So there is no ILO between $`|\mathrm{\Theta }_4`$ and $`|\mathrm{\Theta }_5.`$ Similarly, for the case of $`M=2`$, the same conclusion can be obtained. Q.E.D. |
warning/0506/quant-ph0506037.html | ar5iv | text | # Suppressing decoherence of quantum algorithms by jump codes
## I Introduction
Overcoming decoherence originating from uncontrolled couplings between a quantum system and its environment is one of the major challenges in the realization of quantum computers. Numerous error correcting methods have been designed recently which are able to achieve this goal under ideal circumstances Nielsen and Chuang (2000). In particular, these ideal conditions require that each error has to be corrected immediately and the correction itself has to be performed instantaneously. In practice this requires error detection, syndrome determination, and recovery operation to be executed on a time scale which is vanishingly small compared to the intrinsic time scale of the quantum algorithm which is to be stabilized. Typically it is difficult to fulfil these ideal conditions and so the natural question arises how do error correcting stabilization methods affect quantum algorithms under non-ideal conditions.
In this paper we explore this problem for the recently developed one-error correcting jump codes which have been designed for correcting spontaneous decay processes of qubits Alber et al. (2001). Jump codes are particularly useful in cases in which not only error times but also error positions are known. They are based on an active error correcting quantum code Shor (1995); Knill and Laflamme (1997); Ekert and Jozsa (1996); Knill et al. (2000) which is embedded in a decoherence free subspace Zanardi and Rasetti (1997); Duan and Guo (1997); Lidar et al. (1998) in such a way that all errors taking place between successive spontaneous decay processes are corrected passively. Thus, jump codes require a small number of recovery operations only and, in addition, within the family of all such embedded quantum codes their redundancy is minimal Beth et al. (2003). However, if one was able to control complex many-body Hamiltonians dynamically, it would be possible to correct spontaneous decay processes with quantum codes of even smaller redundancy which involve one redundant qubit only Ahn et al. (2003). Another proposal involving error correction of spontaneous decay processes with one redundant qubit was explored in Ref. Khodjasteh and Lidar (2003). But this suggestion is erroneous as will be discussed later. Thus, as long as it is still difficult to control complicated many-body Hamiltonians jump codes offer attractive perspectives for the correction of spontaneous decay processes as their error correction involves one- and two-qubit Hamiltonians only.
Applying jump codes to the stabilization of quantum algorithms one also ought to be able to correct spontaneous decay processes which occur during the application of elementary quantum gates. For this purpose one has to ensure that even during the application of a quantum gate the error correcting code space is not left at any time Alber et al. (2003a). This requirement can be fulfilled by realizing a universal set of quantum gates with the help of Hamiltonians which leave an error correcting code space invariant. In addition, it is desirable that these Hamiltonians are as simple as possible so that they can be realized in laboratory. Recently, it was demonstrated that such Hamiltonian quantum gates can be constructed in a straightforward way provided one restricts the encoding to appropriate subspaces of one-error correcting jump codes Khodjasteh and Lidar (2002). It is even possible to develop these Hamiltonian universal quantum gates in such a way that the logical qubits constituting these subspaces can be addressed individually, i.e. these logical subspaces can be equipped with a natural tensor-product structure.
In the following we investigate the extent to which these latter one-error correcting jump codes are capable of stabilizing quantum algorithms against spontaneous decay processes under non-ideal conditions. As an example we consider the recently proposed quantum algorithm of the tent-map Frahm et al. (2004). Quantum maps of this kind provide interesting candidates of quantum algorithms which may be run on the first generations of few-qubit quantum computers Benenti et al. (2001, 2003). Even if error times and error positions are known precisely and if the appropriate Hamiltonian quantum gates operate perfectly, residual errors arise due to the finite duration of realistic recovery operations. In particular, any spontaneous decay process occurring during a recovery operation cannot be corrected by an encoding within a one-error correcting jump code. It is demonstrated that the resulting decoherence can be suppressed significantly by using a parallel encoding of the quantum registers of a quantum computer. If the physical qubits of each quantum register constitute a one-error correcting jump code, for example, all simultaneous or sequential spontaneous decay processes can be corrected by such an encoding provided they affect different quantum registers. For this purpose we present a universal entanglement gate which is capable of entangling any logical qubits of any two different one-error correcting jump codes.
This paper is organized as follows: In Sec. II basic aspects of one-error correcting jump codes are summarized and a universal set of quantum gates is discussed whose one- and two-qubit Hamiltonians leave certain subspaces of these jump codes invariant. In addition, a universal entanglement gate is presented which is capable of entangling any two logical qubits belonging to two different error correcting code spaces. In Sec. III the dynamics of the quantum algorithm of the tent-map are explored under the influence of realistic recovery operations of finite duration. It is demonstrated that decoherence can be suppressed significantly by a parallel encoding of the quantum registers which also allows to correct simultaneous spontaneous decay processes affecting different error correcting code spaces.
## II Spontaneous decay of qubits and one-error correcting jump-codes
In this section basic aspects of the recently introduced one-error correcting jump codes are summarized. In particular, the recently proposed logical subspaces Khodjasteh and Lidar (2002) are discussed which can be equipped with a tensor-product structure with the help of universal Hamiltonian quantum gates leaving these subspaces invariant. A novel universal entanglement gate is presented which is capable of entangling arbitrary logical qubits of different logical subspaces. Decomposing quantum algorithms with the help of this quantum gate one can correct simultaneous spontaneous decay processes provided they affect physical qubits of different error correcting code spaces.
Let us consider a typical quantum optical model of a quantum information processor consisting of $`n_q`$ two level quantum systems (qubits) which can decay spontaneously by emission of photons. If the distance between the qubits is much larger than the wave length of the spontaneously emitted photons the resulting spontaneous decay processes are statistically independent. In the Born- and Markov-approximation the dynamics of such a quantum system can be described by a master equation in Lindblad form Carmichael (2003), i.e.
$`\dot{\rho }_t`$ $`=`$ $`{\displaystyle \frac{i}{\mathrm{}}}[H_{sys}(t),\rho _t]+{\displaystyle \underset{\alpha =0}{\overset{n_q1}{}}}([L_\alpha ,\rho _tL_\alpha ^{}]+[L_\alpha \rho _t,L_\alpha ^{}]).`$
(1)
Thereby, $`\rho _t`$ denotes the reduced density operator of the $`n_q`$-qubit quantum system at time $`t`$ and the Hamiltonian $`H_{sys}(t)`$ is assumed to describe the ideal dynamics of the qubits due to a particular quantum algorithm. The Lindblad operator
$$L_\alpha =\sqrt{\kappa _\alpha }|0_{\alpha }^{}{}_{\alpha }{}^{}1|𝟙_{\beta \alpha },$$
(2)
describes the spontaneous decay of qubit $`\alpha `$ from its (unstable) excited state $`|1_\alpha `$ to its (stable) ground state $`|0_\alpha `$ with the spontaneous decay rate $`\kappa _\alpha `$.
For the important special case of equal decay rates, i.e. $`\kappa _\alpha \kappa `$ for $`\alpha =0,\mathrm{},n_q1`$, such a quantum information processor can be protected against spontaneous decay processes with the help of the recently developed one-error detected-jump-error-correcting quantum codes or jump codes (JCs) Alber et al. (2001, 2003b). These quantum codes correct all errors originating from the Lindblad operators of Eq.(1) occurring between successive spontaneous decay processes passively with the help of appropriately constructed decoherence free subspaces (DFSs) Zanardi and Rasetti (1997); Duan and Guo (1997); Lidar et al. (1998). In addition, provided error positions are known these quantum codes are capable of correcting any single spontaneous emission event actively by an appropriately constructed active quantum code which is embedded within a DFS. The orthonormal logical basis states (code words) of these quantum codes are constructed from all possible complementary pairings of the $`n_q`$-qubit states (with $`n_q`$ being even) which involve precisely $`n_q/2`$ excited qubits. In the case of $`n_q=4`$ the orthonormal basis states of the $`(4,3,1)_2`$ code, for example, are given by
$`|c_0_L`$ $`={\displaystyle \frac{1}{\sqrt{2}}}\left(|0011+|1100\right),`$ (3a)
$`|c_1_L`$ $`={\displaystyle \frac{1}{\sqrt{2}}}\left(|0101+|1010\right),`$ (3b)
$`|c_2_L`$ $`={\displaystyle \frac{1}{\sqrt{2}}}\left(|0110+|1001\right).`$ (3c)
Analogously, one may construct one-error correcting $`(n_q,`$ $`\text{dim}_{n_q},1)_{n_q/2}`$-codes for any even number of physical qubits $`n_q`$ with the dimension of the logical Hilbert space
$$\text{dim}_{n_q}=\frac{1}{2}\left(\genfrac{}{}{0pt}{}{n_q}{n_q/2}\right).$$
(4)
This particular family of quantum codes has the interesting property that for any even number of physical qubits $`n_q`$ the redundancy of the associated JC is as small as possible provided one aims at correcting errors between successive spontaneous emission events passively Beth et al. (2003). Furthermore, if a spontaneous decay process of qubit $`\alpha `$ has been detected the resulting error can be corrected by applying the unitary recovery operation
$`R_\alpha =X_\alpha \left({\displaystyle \underset{\beta \alpha }{}}\text{CNot}_{[\alpha ]\beta }\right)\text{H}_\alpha ,`$ (5)
as illustrated in figure 1.
Thereby, $`X_\alpha `$, $`\text{CNot}_{[\alpha ]\beta }`$, and $`\text{H}_\alpha `$ denote a not-gate acting on qubit $`\alpha `$, a controlled not-gate acting on qubit $`\beta `$ with control qubit $`\alpha `$, and a Hadamard-gate acting on qubit $`\alpha `$.
In order to use a one-error correcting $`(n_q,\text{dim}_{n_q},`$ $`1)_{n_q/2}`$-code for the perfect stabilization of arbitrary quantum algorithms one has to develop appropriate sets of universal quantum gates which guarantee that the code space is not left even during the application of these gates Alber et al. (2003a). Otherwise any spontaneous emission process taking place during the application of an elementary quantum gate can no longer be corrected. This requirement can be achieved with the help of Hamiltonian quantum gates provided the relevant Hamiltonians leave the relevant code space invariant. So far, for general one-error correcting jump codes the construction of such universal quantum gates has been possible in special cases only Alber et al. (2003a). However, restricting oneself to appropriate subspaces of these one-error correcting jump codes it is possible to construct such universal quantum gates in a straightforward way at the prize of increasing redundancy. Examples of such subspaces have been proposed recently by Khodjasteh and Lidar Khodjasteh and Lidar (2002). In the case of $`n_q`$ physical qubits such an appropriate subspace of the one-error correcting $`(n_q,\text{dim}_{n_q},1)_{n_q/2}`$-code is spanned by the orthonormal logical states
$$\begin{array}{cc}\hfill |00\mathrm{}0_L& =\frac{1}{\sqrt{2}}(|01|01\mathrm{}|01|01+\hfill \\ & |10|10\mathrm{}|10|10),\hfill \\ \hfill |00\mathrm{}1_L& =\frac{1}{\sqrt{2}}(|01|01\mathrm{}|10|01+\hfill \\ & |10|10\mathrm{}|01|10),\hfill \\ & \mathrm{}\hfill \\ \hfill |11\mathrm{}1_L& =\frac{1}{\sqrt{2}}(|10|10\mathrm{}|10|01+\hfill \\ \hfill \underset{\begin{array}{c}n_L\text{ logical}\\ \text{qubits}\end{array}}{\underset{}{}}& \underset{2n_L\text{ physical qubits}}{\underset{}{|01|01\mathrm{}|01}}|10)\hfill \end{array}$$
(6)
for example. In this case $`n_q`$ physical qubits are required for the encoding of $`n_L=(n_q2)/2`$ logical qubits. Besides the complementary pairing of the $`n_q`$-qubit states Eq.(6) also involves a complementary pairing of adjacent physical qubits. Thereby the rightmost two physical qubits are used for distinguishing the complementary states of each pair of $`n_q`$-qubit states. In the code space spanned by the orthonormal logical states of Eq.(6) any spontaneous decay process of any of the $`n_q`$ physical qubits can be corrected provided error position and error time are known.
### II.1 Universal sets of quantum gates within one-error correcting code spaces
How can we develop a Hamiltonian set of universal quantum gates which do not leave the code space of Eq.(6) at any time? This can be achieved with the help of the elementary Ising- and Heisenberg-type two-qubit Hamiltonians $`T_{\alpha \beta }=(X_\alpha X_\beta +Y_\alpha Y_\beta )/2`$ and $`ZZ_{\alpha \beta }Z_\alpha Z_\beta `$, for example, which act on the physical qubits $`\alpha `$ and $`\beta `$. (The operators $`X,Z`$, and $`Y`$ denote the three anti-commuting Pauli spin operators of the appropriate qubits with $`XY=iZ`$ etc.) Provided one is capable of controlling these two-qubit Pauli Hamiltonians it is straightforward to construct a universal set of logical single and two-qubit Hamiltonians $`\overline{X}_i,\overline{Z}_i`$, and $`\overline{ZZ}_{ij}`$ which act on the logical qubits $`i`$ and $`j`$ similarly as the corresponding Hamiltonians $`T_{\alpha \beta },ZZ_{\alpha \beta }`$ and which leave the code space of $`n_L`$ logical qubits invariant Khodjasteh and Lidar (2002), i.e.
$`\overline{X}_i`$ $`=T_{2i+3,2i+2},`$ (7a)
$`\overline{Z}_i`$ $`=ZZ_{2i+3,1},`$ (7b)
$`\overline{ZZ}_{ij}`$ $`=Z_{2i+3}Z_{2j+3}.`$ (7c)
With the help of these Hamiltonians any two logical qubits $`i`$ and $`j`$ can be addressed individually thus providing a tensor product structure in the $`2^{n_L}`$-dimensional Hilbert space spanned by the basis state of Eq.(6). In Eq.(6) it is assumed that the logical (physical) qubits are numbered from right to left from $`i,j=0`$ ($`\alpha ,\beta =0`$) up to the maximum value of $`i,j=n_L1`$ ($`\alpha ,\beta =n_q1`$). The names of these Hamiltonians indicate how they act on the logical states. Thus, the Hamiltonian $`\overline{ZZ}_{ij}`$, for example, acts on the logical qubits $`i`$ and $`j`$ in the same way as the Hamiltonian $`ZZ_{\alpha \beta }`$ acts on the physical qubits $`\alpha `$ and $`\beta `$.
Quantum gates built with the help of these Hamiltonians always stay within the code space spanned by the logical basis of Eq.(6) for any even number of physical qubits. The construction of universal logical single and two-qubit gates, such as Hadamard- and phase gates, based on these Hamiltonians is exemplified in Figs. 2 and 3.
Furthermore, as the code space of Eq.(6) is a subspace of a $`(n_q,\mathrm{dim}_{n_q},1)_{n_q/2}`$ one-error correcting jump code any detected spontaneous decay of any physical qubit $`\alpha `$ can be corrected by the recovery operation of Eq.(5). On this occasion we want to point out that the recovery operation proposed in Ref. Khodjasteh and Lidar (2002) does not work properly since it does not even restore the basis states of Eq.(6). This mistake might have its origin in a misinterpretation of the jump operator $`|0`$$`1|`$ which actually transforms the physical state $`|0`$ into the zero-vector $`0`$ (which is not a physical state) and not into $`|0`$ again. The same misinterpretation appears to lead to the construction of a wrong code space in a subsequent paper by the same authors Khodjasteh and Lidar (2003).
### II.2 A universal entanglement gate for different one-error correcting code spaces
The one-error correcting quantum codes of Eq.(6) cannot correct simultaneous spontaneous decay processes affecting different physical qubits. However, this problem can be overcome at least partly by combining one-error correcting quantum codes each of which involves a relatively small number of physical qubits. Physically this can be achieved by a local architecture of a quantum information processor, for example, which is based on small (local) groups of physical qubits (elementary registers) each of which constitutes a one-error correcting code space as described by Eq.(6). As a result any number of simultaneous decay processes can be corrected in parallel provided these decays affect physical qubits of different one-error correcting code spaces. In order to be able to stabilize arbitrary quantum algorithms in such a quantum information processor one has not only to be able to perform arbitrary unitary operations within each elementary quantum register but one also has to be able to entangle any two logical qubits of any to different elementary quantum registers. In the following we construct such a universal entanglement gate.
For this purpose let us consider first of all the simple case of the tensor product space of two logical qubits. According to Eq.(6) the resulting four basis states are given by
$$\begin{array}{cc}\hfill |0_{La}|0_{Lb}& =\frac{1}{2}(|01|01+|10|10)_a\hfill \\ & \left(|01|01+|10|10\right)_b,\hfill \\ \hfill |0_{La}|1_{Lb}& =\frac{1}{2}(|01|01+|10|10)_a\hfill \\ & \left(|10|01+|01|10\right)_b,\hfill \\ \hfill |1_{La}|0_{Lb}& =\frac{1}{2}(|10|01+|01|10)_a\hfill \\ & \left(|01|01+|10|10\right)_b,\hfill \\ \hfill |1_{La}|1_{Lb}& =\frac{1}{2}(|10|01+|01|10)_a\hfill \\ & \left(|10|01+|01|10\right)_b\hfill \end{array}$$
(8)
with the subscripts $`a`$ and $`b`$ referring to two different quantum registers. The Hamiltonian
$$H_{\text{ent}}^{a,b}=ZZ_{4,0}+ZZ_{6,1}+ZZ_{5,2}+ZZ_{7,3}$$
(9)
has the following properties
$`H_{\text{ent}}^{a,b}`$ $`|m_{La}|n_{Lb}=0,(m,n)(1,1)`$
$`H_{\text{ent}}^{a,b}`$ $`|u_{a,b}=4|u_{a,b},`$
$`H_{\text{ent}}^{a,b}`$ $`|v_{a,b}=4|v_{a,b}`$
with $`|u_{a,b}=(|10|01)_a(|10|01)_b+(|01|10)_a(|01|10)_b`$, $`|v_{a,b}=(|10|01)_a(|01|10)_b+(|01|10)_a(|10|01)_b`$. Thus, this Hamiltonian generates a controlled $`\pi `$-phase gate
$$\text{CP}\left(\pi \right)𝟙\mathrm{𝟚}|1111|=\mathrm{exp}\left(𝕚_{\text{ent}}^{𝕒,𝕓}\frac{\pi }{\mathrm{𝟜}}\right).$$
(11)
Since the four $`ZZ`$ Hamiltonians appearing in Eq.(9) commute, they can also be applied to the physical qubits one after the other. Note that the gate of Eq.(11) has similarities with the entanglement gate presented in Ref. Alber et al. (2003a) for entangling the qubits of two one-error correcting $`(4,3,1)_2`$-code spaces.
The Hamiltonian entanglement gate of Eq.(11) can be generalized to two jump codes of Eq.(6) of arbitrary sizes. In particular, a controlled $`\pi `$-phase gate affecting the logical qubits $`j_a`$ and $`j_b`$ of the elementary registers $`a`$ and $`b`$ which contain $`n_L`$ and $`\stackrel{~}{n}_L`$ logical qubits is given by Eq.(11) with the entangling Hamiltonian (compare with Fig.4)
$`H_{\text{ent}}^{j_a,j_b}`$ $`=`$ $`ZZ_{2\stackrel{~}{n}_L+2,0}+ZZ_{2j_a+2\stackrel{~}{n}_L+4,1}+`$ (12)
$`ZZ_{2\stackrel{~}{n}_L+3,2j_b+2}+ZZ_{2j_a+2\stackrel{~}{n}_L+5,2j_b+3}.`$
Together with the universal quantum gates based on the Hamiltonians of Eqs.(7a),(7b),(7c) this controlled $`\pi `$-phase gate constitutes a universal set of elementary quantum gates Brylinski (2002) with which any unitary transformation can be realized in a quantum information processor with a local architecture. Two examples of such unitary transformations and their corresponding decompositions in terms of these universal quantum gates are depicted in Fig. 5.
There is one particularity concerning the recovery operation which must be applied to restore the code space after a detected spontaneous decay taking place during the application of the entanglement gate of Eq.(11). In general during the application of this entanglement gate a quantum state leaves the tensor product space of the associated elementary quantum registers. But a perfect recovery by an appropriate gate sequence (5) is still possible. This is due to the fact that even during the application of an entanglement gate of the form of Eq.(12) an arbitrary linear superposition of basis states of Eq.(6) always remains inside the code space of a one-error correcting jump code.
## III Stabilization of quantum maps under non-ideal conditions
The jump codes discussed in the previous section provide perfect protection against spontaneous decays of qubits provided error times and error positions are known precisely and appropriate recovery operations are applied immediately and instantaneously. In this section we investigate the dynamics of quantum algorithms whose decoherence with respect to spontaneous decay processes is stabilized by the one-error correcting jump codes discussed in Sec. II. In particular, we are mainly interested in situations in which a quantum algorithm is implemented ideally by a universal set of quantum gates which do not leave an error correcting code space at any time. This way one is able to correct spontaneous decay processes even if they occur during the application of an elementary quantum gate. However, as a result of such an encoding the time required for the application of a recovery operation is typically no longer negligibly small in comparison with the intrinsic time evolution of a quantum algorithm so that the ideal conditions of error correction are no longer fulfilled.
In order to explore the consequences of such non-ideal conditions let us consider the recently proposed quantum algorithm of the iterated tent-map Frahm et al. (2004) as an example. In each iteration of the quantum tent-map an initial quantum state $`|\psi `$ is mapped onto the state
$`|\psi ^{}`$ $`=`$ $`e^{iT\widehat{p}^2/2}e^{ik\widehat{V}(\widehat{x})}|\psi `$ (13)
with the tent-shaped force
$`V^{}(x)=\{\begin{array}{cc}(x\frac{\pi }{2}),\hfill & (0x<\pi )\hfill \\ (\frac{3\pi }{2}x),\hfill & (\pi x<2\pi ).\hfill \end{array}`$ (16)
Thereby, $`\widehat{p}`$ and $`\widehat{x}`$ are dimensionless quantized action-angle variables. Universal gate sequences involving Hadamard-, phase-, controlled phase-, and controlled not-gates were developed recently Frahm et al. (2004). In order to stabilize this quantum algorithm against spontaneous decay processes these universal quantum gates and the recovery operations required have to be decomposed in terms of the elementary gates discussed in the previous section which do not leave the error correcting code space at any time (compare with Figs.2, 3, 5).
In the following we present numerical simulations for $`t=30`$ iterations of the quantum tent-map involving $`n_L=6`$ logical qubits. The characteristic parameters chosen are $`k=1.7/T`$ and $`T=2\pi /2^{n_L}`$. As an initial quantum state we chose a coherent state centered around the mean value $`(x,p)=(5.35,0)`$ Frahm et al. (2004). We compare the dynamics of the quantum algorithm which is stabilized by the one-error correcting jump code of Eq.(6) involving $`n_q=2n_L+2=14`$ physical qubits with the corresponding dynamics of $`n_q=n_L=6`$ qubits in the absence of error correction for different spontaneous decay rates $`\kappa `$. With this choice of parameters each iteration of the quantum tent-map can be decomposed into $`n_g=125`$ elementary Hadamard-, phase-, controlled phase-, and controlled not-gates acting on these $`n_L`$ logical qubits Frahm et al. (2004). In our numerical simulations each of these $`n_g`$ quantum gates is decomposed into a suitable sequence of universal logical quantum gates according to the gate libraries depicted in Figs. 2 and 3. Thus, for each iteration of the quantum tent-map $`437`$ of these universal quantum gates are required whose Hamiltonians have to be turned on for appropriate values of the dimensionless parameter $`\tau `$. This latter parameter may be viewed as a dimensionless measure for the duration of each of these elementary universal quantum gates. Adding the $`\tau `$-values of the gates needed for one iteration of the quantum map one obtains the value $`\tau _{it}=67.2\pi `$ which implies an average dimensionless time of magnitude $`\tau _{it}/n_g=0.54\pi `$ for each of the $`n_g`$ quantum gates of Ref. Frahm et al. (2004).
The numerical results are obtained by simulating the master equation of Eq.(1) with the quantum trajectory method Dum et al. (1992). As soon as a spontaneous decay process takes place the relevant Hamiltonian, say $`\widehat{H}_0`$, implementing the actual logical quantum gate is stopped immediately and the Hamiltonians are turned on which are required to perform the relevant recovery operation. Summing up the $`\tau `$-values of a recovery operation yields $`\tau _{\text{rec}}=n_q\tau _{\text{cnot}}\pi /2=24\pi `$ (for $`n_q=14`$) with the duration of a CNOT-operation being given by $`\tau _{\text{cnot}}=7\pi /4`$ (compare with Eq. (5) and Figs. 2 and 3). After completion of the recovery operation the quantum algorithm is started again with the previously stopped Hamiltonian $`\widehat{H}_0`$. Of course, spontaneous decay processes occurring during a recovery operation cannot be corrected. The resulting non-ideal performance of the quantum algorithm can be measured by the fidelity
$$f(t)=\mathrm{Tr}(\rho _t|\overline{\psi }(t)\overline{\psi }(t)|)$$
(17)
of the non-ideal quantum state $`\rho _t`$ with respect to the corresponding ideal (pure) quantum state $`|\overline{\psi }(t)`$.
In Fig. 6 the time evolution of this fidelity is depicted for different values of the (dimensionless) spontaneous decay rate $`\kappa `$ of the qubits without (upper diagram) and with (lower diagram) error correction. The fidelities presented were averaged over $`10^3`$ statistical realizations. In the upper diagram of Fig. 6 $`|\overline{\psi }(t)`$ denotes the ideal quantum state of a quantum computer consisting of $`n_q=6`$ physical qubits and $`\rho _t`$ is the corresponding quantum state in the presence of spontaneous emission. In the lower diagram of Fig.6 $`|\overline{\psi }(t)`$ is the ideal quantum state of a quantum computer consisting of $`n_q=14`$ physical qubits in the absence of spontaneous decay processes. Analogously, $`\rho _t`$ is the quantum state of a quantum computer involving $`n_q=14`$ physical and $`n_L=(n_q2)/2=6`$ logical qubits whose dynamics are stabilized against spontaneous decay processes by an appropriate encoding with the help of the one-error correcting quantum code of Eq.(6) and the quantum gates of Figs. 2 and 3.
It is apparent from the dashed curves in Fig. 6 that the numerically obtained fidelities are described approximately by the relations
$$f_{\text{er}}(t)=\mathrm{exp}\left(\frac{n_q}{2}\kappa \tau _{\text{it}}t\right)$$
(18)
without error correction and by
$$f_{\text{ec1}}(t)=\mathrm{exp}\left(\left(\frac{n_q}{2}\kappa \right)^2\tau _{\text{rec}}\tau _{\text{it}}t\right)e^{R_{\text{ec1}}\tau _{\text{it}}t},$$
(19)
in the case of error correction. Thereby, $`t`$ denotes the number of iterations of the quantum map and $`R_{\text{ec1}}`$ is the fidelity decay rate in the presence of error correction. The quantity $`f_{\text{er}}(t)`$ of Eq.(18) resembles the recently proposed formula of Ref. Lee and Shepelyansky (2005) where the elementary quantum gates were performed instantaneously. In our notation this corresponds to the case $`\tau _{\text{it}}=n_g`$. According to Eq.(18) the mean fidelity is estimated by the probability that no spontaneous decay process takes place in a time interval of duration $`\tau _{\text{it}}\times t`$ for a quantum state with $`n_q/2`$ excited qubits. This estimate is based on the assumption that a fidelity of unity is associated with those particular statistical realizations for which no spontaneous decay takes place during the time interval $`\tau _{\text{it}}\times t`$. All other statistical realizations are assumed to yield a zero fidelity. The estimate for $`f_{\text{ec1}}(t)`$ of Eq.(19) is based on the analogous assumption that a unit fidelity is associated with those statistical realizations only in which no spontaneous decay process takes place during a recovery operation requiring a time $`\tau _{\text{rec}}`$. If only one spontaneous decay process were possible, the associated probability would be given by $`p=\mathrm{exp}(\kappa \tau _{\text{rec}}n_q/2)`$ for $`n_q/2`$ excited qubits. However, on the average there are $`N=\kappa \tau _{\text{it}}tn_q/2`$ spontaneous decay processes taking place during the time interval of interest. Provided these decays are statistically independent the mean fidelity can be estimated by $`p^N`$ which yields Eq.(19). From the approximate relations (18) and (19) it is apparent that the use of the one-error correcting code space of Eq.(6) together with the quantum gates of Figs.2 and 3 implies a significant increase of the fidelity as long as the average time between two successive spontaneous decay processes is much larger than the recovery time $`\tau _{\text{rec}}=n_q\tau _{\text{cnot}}\pi /2`$. For the duration $`\tau _{\text{cnot}}`$ of the CNOT-operations involved in a recovery operation this condition yields
$$\tau _{cnot}\frac{1}{4\kappa }\frac{n_L}{(n_L+1)^3}.$$
(20)
Thus, for large numbers of logical qubits $`n_L`$ it may be difficult to realize such fast CNOT-operations.
Relation (19) is applicable if spontaneous decay processes occurring during recovery operations are the main source of decoherence. Using an encoding based on groups of physical qubits, i.e. quantum registers, each of which involves its own one-error correcting code space at least some of the spontaneous decay processes occurring during recovery operations can be corrected provided they take place in different registers. This way an additional significant improvement of the stabilization of quantum algorithms may be obtained if spontaneous decay during recovery operations is the main reason for decoherence. For this purpose let us consider $`N_{\text{reg}}`$ basic quantum registers each of which consists of $`n_{\text{reg}}^{(i)}`$ physical qubits and encodes $`n_L^{(i)}=(n_{\text{reg}}^{(i)}2)/2`$ logical qubits according to Eq.(6). Local operations within each quantum register can be performed with the help of the universal quantum gates of Figs. 2 and 3. Entanglement operations between different quantum registers can be performed with the entanglement gate of Fig.5. In such an architecture of a quantum information processor all spontaneous decay processes occurring in different quantum registers can be corrected irrespective of whether they occur successively or simultaneously. According to the reasoning used for the derivation of Eq. (19) the mean fidelity decay of such an encoding can be estimated by
$`f(t)`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{N_{\text{reg}}}{}}}(e^{(\kappa /2)n_{\text{reg}}^{(i)}\tau _{\text{rec}}^{(i)}})^{N^{(i)}}e^{R_{\text{ec2}}\tau _{\text{it}}t}`$
with the fidelity decay rate
$`R_{\text{ec2}}`$ $`=`$ $`(\kappa /2)^2{\displaystyle \underset{i=1}{\overset{N_{\text{reg}}}{}}}\tau _{\text{rec}}^{(i)}(n_{\text{reg}}^{(i)})^2,`$ (22)
the recovery time $`\tau _{\text{rec}}^{(i)}=n_{\text{reg}}^{(i)}\tau _{\text{cnot}}\pi /2`$ and with $`N^{(i)}=(\kappa /2)n_{\text{reg}}^{(i)}\tau _{\text{it}}t`$ denoting the mean number of spontaneous decay processes taking place within register $`i`$ during $`t`$ iterations of the quantum map. If all quantum registers involve the same total number of logical qubits $`n_L^{(i)}`$, according to Eqs.(18) and (LABEL:block) a blockwise encoding implies a significant increase of fidelity as long as the duration of the CNOT-operations involved in the recovery operations are small enough, i.e.
$`\tau _{\text{cnot}}{\displaystyle \frac{1}{4\kappa }}{\displaystyle \frac{n_L^{(i)}}{(n_L^{(i)}+1)^3}}.`$ (23)
Thus, contrary to the rather stringent condition (20) for $`n_L^{(i)}=1`$, for example, successful error correction by blocks of one-error correcting jump codes requires only that the duration of a CNOT-operation $`\tau _{\text{cnot}}`$ is much smaller than the mean life time of a single qubit $`1/\kappa `$. In realistic applications it should not be too difficult to fulfil this condition. So, despite its higher redundancy for any fixed total number of logical qubits $`n_L`$ a blockwise encoding based on $`n_L/n_L^{(i)}`$ quantum registers each of which contains $`n_L^{(i)}n_L`$ logical qubits is more stable against spontaneous decay processes than a direct encoding of these $`n_L`$ logical qubits with $`n_q=2n_L+2`$ physical qubits according to Eq.(6).
## IV Summary and conclusions
One-error correction jump codes are an efficient method to stabilize quantum algorithms against spontaneous decay processes provided one is capable of correcting also errors taking place during elementary quantum gates. With the help of appropriate subspaces of these error correcting code spaces simple Hamiltonian universal quantum gates can be constructed which achieve this goal. A realistic simulation was presented in which the quantum algorithm of the tent-map was decomposed into these quantum gates. Even if the recovery operations cannot be implemented ideally this encoding suppresses the decohering influence of spontaneous decay processes significantly. This error suppression can be increased with the help of a block-encoding which involves different error correcting code spaces. For this purpose the presented universal entanglement gate might turn out to be useful as it allows to suppress simultaneous spontaneous decay processes provided they affect physical qubits of different error-correcting logical subspaces.
###### Acknowledgements.
This work is supported by the EU IST-FET project EDIQIP and by the DFG (SPP-QIV). Numerical support by Stanislav Vymětal and Igor Jex is gratefully acknowledged. |
warning/0506/math0506587.html | ar5iv | text | # 1 Non-parametric settings
## 1 Non-parametric settings
For natural $`n3,`$ let us consider a two-dimensional minimal surface graph
$$X(x,y)=(x,y,\phi _1(x,y),\mathrm{},\phi _{n2}(x,y)),(x,y)\overline{B}_R,$$
with functions $`\phi _\mathrm{\Sigma }C^2(\overline{B}_R,),`$ $`\mathrm{\Sigma }=1,\mathrm{},n2,`$ on the topological closure $`\overline{B}_R^2`$ of the open disc
$$B_R:=\{(x,y)^2:x^2+y^2<R^2\},R(0,+\mathrm{}).$$
With the minimal surface we associate the moving frame $`\{X_x,X_y,N_1,\mathrm{},N_{n2}\}`$ consisting of the tangent vectors
$$X_x=(\mathrm{1,0},\phi _{1,x},\mathrm{},\phi _{n2,x}),X_y=(\mathrm{0,1},\phi _{1,y},\mathrm{},\phi _{n2,y}),$$
where the indices $`x`$ and $`y`$ denote the partial derivatives w.r.t. the coordinates, and the linearly independent unit normal vectors
$$\begin{array}{ccc}\hfill N_1& :=& \frac{1}{\sqrt{1+|\phi _1|^2}}(\phi _{1,x},\phi _{1,y}\mathrm{,1,0},\mathrm{}\mathrm{,0,0}),\hfill \\ \hfill \mathrm{}& & \mathrm{}\hfill \\ \hfill N_{n2}& :=& \frac{1}{\sqrt{1+|\phi _{n2}|^2}}(\phi _{n2,x},\phi _{n2,y}\mathrm{,0,0},\mathrm{}\mathrm{,0,1}).\hfill \end{array}$$
(1)
## 2 Introduction of conformal parameters
Instead of the coordinates $`x`$ and $`y,`$ we prefer to consider the graph conformally parametrized (see ). Let $`(u,v)\overline{B},`$ where $`B:=B_1,`$ then we get the immersion
$$X(u,v)=(x^1(u,v),x^2(u,v),x^3(u,v),\mathrm{},x^n(u,v)),(u,v)\overline{B},$$
of regularity class $`XC^2(B,^n)C^0(\overline{B},^n),`$ while the following properties hold true:
* the plane mapping
$$F^{}(u,v):=(x^1(u,v),x^2(u,v)):\overline{B}\overline{B}_R$$
is one-to-one, it has a positive Jacobian, it satisfies $`F^{}(\mathrm{0,0})=(\mathrm{0,0}),`$ and $`F^{}:BB_R`$ is a topological mapping, where
$$B_R:=\{(x,y)^2:x^2+y^2=R^2\}$$
denotes the boundary of $`B_R;`$
* the conformality relations
$$|X_u|^2=W=|X_v|^2,X_uX_v^t=0\text{in}B$$
with the surface area element
$$W:=\sqrt{h_{11}h_{22}h_{12}^2},h_{ij}:=X_{u^i}X_{u^j}^t\text{for}i,j=\mathrm{1,2}$$
and $`u^1u,`$ $`u^2v.`$
Now, consider any differentiable unit normal vector field $`N=N(u,v).`$ W.r.t. its direction we define the coefficients of the second fundamental form by
$$L_{N,ij}:=X_{u^iu^j}N^t,i,j=\mathrm{1,2}.$$
The mean curvature and the Gaussian curvature in direction of $`N`$ read
$$H_N:=\frac{L_{N\mathrm{,11}}h_{22}2L_{N\mathrm{,12}}h_{12}+L_{N\mathrm{,22}}h_{11}}{2W^2},K_N:=\frac{L_{N\mathrm{,11}}L_{N\mathrm{,22}}L_{N\mathrm{,12}}^2}{W^2}.$$
In terms of the corresponding principle curvatures $`\kappa _{N\mathrm{,1}}`$ and $`\kappa _{N\mathrm{,2}},`$ we can write
$$H_N=\frac{\kappa _{N\mathrm{,1}}+\kappa _{N\mathrm{,2}}}{2},K_N=\kappa _{N\mathrm{,1}}\kappa _{N\mathrm{,2}}.$$
###### Lemma.
Let $`XC^2(B,^n)C^0(\overline{B},^n)`$ be a minimal immersion. Then there hold
$$H_N(u,v)0\text{in}B$$
for all differentiable unit normal vector fields $`N=N(u,v).`$
For the proof of this Lemma and further results on two-dimensional minimal immersions we refer to the fundamental paper .
## 3 The main result
The main object of this note is the following
###### Theorem.
Let the minimal surface graph $`X(x,y)=(x,y,\phi _1(x,y),\mathrm{},\phi _{n2}(x,y))`$ with $`\phi _\mathrm{\Sigma }C^2(\overline{B}_R,),`$ $`\mathrm{\Sigma }=1,\mathrm{},n2,`$ be given. Then, for the principle curvatures $`\kappa _{N\mathrm{,1}}`$ and $`\kappa _{N\mathrm{,2}}`$ in direction of any unit normal vector field $`N,`$ there is a universal constant $`\mathrm{\Theta }(0,+\mathrm{})`$ such that in terms of conformal parameters $`(u,v)\overline{B}`$ it holds
$$\kappa _{N\mathrm{,1}}(\mathrm{0,0})^2+\kappa _{N\mathrm{,2}}(\mathrm{0,0})^2\frac{1}{R^4}\mathrm{\Theta }X_{C^0(B)}^2$$
with the Schauder norm
$$X_{C^0(B)}:=\underset{(x,y)B}{sup}|X(x,y)|.$$
We conclude the following result of Bernstein-Liouville type:
###### Corollary.
For large $`R,`$ let $`X_{C^0(B_R)}\mathrm{\Omega }R^\omega `$ with a real constant $`\mathrm{\Omega }(0,+\mathrm{})`$ and $`\omega [\mathrm{0,2}).`$ Then, if the minimal graph is complete, that is it is defined over the whole plane $`^2,`$ it represents an affine plane.
First, we continue with some remarks.
###### Remarks.
* The corollary is sharp in the following sense: the holomorphic function $`zz^2,`$ $`z=x+iy,`$ induces the minimal graph $`(z,z^2)`$ in $`^4`$ which is not linear.
* Graphs $`(z,\mathrm{\Phi }(z)),`$ where $`\mathrm{\Phi }=\mathrm{\Phi }(z)`$ is any holomorphic function, are minimal in $`^4.`$ Thus, the above corollary can be read as a generalization of the well-known Liouville theorem.
* In , using methods from complex function theory it was proved a curvature estimate and a theorem of Bernstein type for two-dimensional minimal immersions in $`^n`$ under the assumption that all normals of such a surface omit any given neighborhood of some direction in space.
In , Osserman’s condition enabled us to establish upper curvature bounds for immersions of prescribed mean curvature fields. The assumptions of our theorem above are more restrictive but the proof of the theorem seems to be simpler.
###### Proof of the Theorem.
* Note that
$$\begin{array}{ccc}|K_N(\mathrm{0,0})|\hfill & =\hfill & \frac{|L_{N\mathrm{,11}}(\mathrm{0,0})L_{N\mathrm{,22}}(\mathrm{0,0})L_{N\mathrm{,12}}(\mathrm{0,0})^2|}{W(\mathrm{0,0})^2}\hfill \\ & \hfill & \frac{|X_{uu}(\mathrm{0,0})||X_{vv}(\mathrm{0,0})|+|X_{uv}(\mathrm{0,0})|^2}{W(\mathrm{0,0})^2}.\hfill \end{array}$$
Thus, we have to establish an upper bound for the second derivatives $`X_{uu}(\mathrm{0,0})`$ etc., and a lower bound for $`W(\mathrm{0,0}).`$
First, a minimal surface in conformal parameters is harmonic, that is (see )
$$\mathrm{}X(u,v)=0\text{in}B.$$
Applying , Theorem 4.6, there is a constant $`C_1(0,+\mathrm{})`$ such that
$$|X_{u^iu^j}(\mathrm{0,0})|C_1X_{C^0(B)},i,j=\mathrm{1,2}.$$
* Instead of the plane mapping $`F^{}=F^{}(u,v)`$ we consider the normalization
$$F(u,v):=\frac{1}{R}F^{}(u,v),(u,v)\overline{B}.$$
Together with the properties of $`F^{}`$ mentioned above in point (a), due to , vol. II, chapter XII, Satz 1 there is a universal constant $`C_2(0,+\mathrm{})`$ such that
$$|F(\mathrm{0,0})|^2C_2$$
(E. Heinz, 1952). Now, using the conformality relations we estimate as follows:
$$\begin{array}{ccc}2W(\mathrm{0,0})\hfill & =\hfill & |x^1(\mathrm{0,0})|^2+|x^2(\mathrm{0,0})|^2+|x^3(\mathrm{0,0})|^2+|x^4(\mathrm{0,0})|^2\hfill \\ & \hfill & |F^{}(\mathrm{0,0})|^2=R^2|F(\mathrm{0,0})|^2R^2C_2.\hfill \end{array}$$
* We arrive at the estimate
$$|K_N(\mathrm{0,0})|\frac{8C_1^2X_{C^0(B)}^2}{R^4C_2^2}.$$
Finally,
$$\kappa _{N\mathrm{,1}}^2+\kappa _{N\mathrm{,2}}^2=2(2H_N^2K_N)=2(K_N)0\text{for any}N$$
yields
$$\kappa _{N\mathrm{,1}}(\mathrm{0,0})^2+\kappa _{N\mathrm{,2}}(\mathrm{0,0})^2\frac{1}{R^4}\frac{16C_1^2}{C_2^2}X_{C^0(B)}^2.$$
Setting $`\mathrm{\Theta }:=16C_1^2C_2^2`$ proves the theorem.
###### Proof of the Corollary.
We introduce conformal parameters $`(u,v)\overline{B}.`$
* Orthonormalization of the unit normal basis given in (1) yields a unit normal frame $`\{N_1(u,v),\mathrm{},N_{n2}(u,v)\}`$ such that
$$N_\mathrm{\Sigma }(u,v)N_\mathrm{\Omega }(u,v)^t=\delta _{\mathrm{\Sigma }\mathrm{\Omega }}:=\{\begin{array}{c}1\text{for}\mathrm{\Sigma }=\mathrm{\Omega }\hfill \\ 0\text{for}\mathrm{\Sigma }\mathrm{\Omega }\hfill \end{array}.$$
The Gaussian curvature $`K`$ (that is the non-trivial component of the Riemann curvature tensor) can be calculated by
$$K=\underset{\mathrm{\Sigma }=1}{\overset{n2}{}}K_\mathrm{\Sigma }$$
with the Gauss curvatures $`K_\mathrm{\Sigma }`$ in direction $`N_\mathrm{\Sigma }`$ (see , equation (6.3,17)). For our minimal surface graph there hold $`H_\mathrm{\Sigma }=0`$ and $`K_\mathrm{\Sigma }0`$ for $`\mathrm{\Sigma }=1,\mathrm{},n2.`$
* For large $`R,`$ we have by assumption $`X_{C^0(B)}\mathrm{\Omega }R^\omega .`$ From the curvature estimate above we conclude
$$\kappa _{N\mathrm{,1}}(\mathrm{0,0})^2+\kappa _{N\mathrm{,2}}(\mathrm{0,0})^2\frac{1}{R^4}R^{2\omega }\mathrm{\Theta }\mathrm{\Omega }^2=\frac{1}{R^{42\omega }}\mathrm{\Theta }\mathrm{\Omega }^2$$
for any $`N=N(u,v)`$ of the orthonormal section $`\{N_1,\mathrm{},N_{n2}\}.`$ Letting $`R\mathrm{}`$ implies $`K_N(\mathrm{0,0})=0`$ and therefore $`K(\mathrm{0,0})=0`$ for the Gaussian curvature. By , chapter III, §5, Lemma 3, the graph is linear.
Steffen Fröhlich
AG Differentialgeometrie und Geometrische Datenverarbeitung
TU Darmstadt
Schloßgartenstraße 7
D-64289 Darmstadt
e-mail: sfroehlich@mathematik.tu-darmstadt.de |
warning/0506/cond-mat0506185.html | ar5iv | text | # Vectorial Loading of Processive Motor Proteins: Implementing a Landscape Picture
## 1 Introduction
A processive motor protein how is an individual protein molecule that in an appropriate aqueous solution binds to a linear periodic directed molecular track and, when fueled via diffusion by suitable molecules (specifically ATP in cases of most interest), takes tens to hundreds of discrete steps along the track before dissociating. Such a motor steps overwhelmingly in a single characteristic direction, which we will identify as parallel to the positive $`x`$-axis: see Figure 1.
Thus conventional kinesin walks towards the plus or fast-growing end of a microtubule while myosin V moves towards the plus or ‘barbed’ end of an actin filament how . In stepping along its track a motor exerts a tension $`𝑭=(F_x,F_y,F_z)`$ on a molecular tether the other end of which is, in vivo, bound to some cellular organelle or vesicle while for in vitro experiments it is firmly attached to a silica or plastic bead that may be controlled with the aid of an optical trap how . For most purposes one may regard the molecular track as fixed in space: relatively rigid microtubules are typically clamped to a microscope slide while the more flexible actin filaments can be stretched between two further beads that are attached to the filament ends and held in place by individual optical traps how .
In the conceptually most straightforward experiments how ; koj:mut ; vis:sch a motor (M in Fig. 1) advances along the track in a positive direction in a solution of fixed fuel concentration, say \[ATP\], under a resistive load the $`x`$-component of which, $`F_x<0`$, is observed or controlled. Measurements of the mean velocity and the dispersion vis:sch ; sch:vis :
$$Vx(t)/t,\text{and}D[x^2(t)x(t)^2]/t,$$
(1)
as functions of $`F_x`$ and \[ATP\], where $`x(t)`$ represents the displacement of the motor along the track at time $`t`$, may then be analyzed fis:kol ; fis:kol2 ; kol:fis with the aim of extracting details of the mechanism by which the motor takes individual steps. In particular one would like to identify and quantify any substeps or intermediate motions.
To investigate further experimentally it is rather natural to impose an oppositely directed reverse or assisting load on the motor so that $`F_x>0`$. Indeed, such experiments have been performed on kinesin, first by Coppin et al. cop:pie , and, more recently, by Block and coworkers lan:asb ; blo:asb , and by Carter and Cross car:cro . In a preliminary computation to gain theoretical insight regarding the effects of reverse loading, Fisher and Kolomeisky (FK) fis:kol2 examined the fairly noisy data of Coppin et al. cop:pie for the mean velocity of kinesin under assisting loads up to $`F_x=+\mathrm{\hspace{0.17em}6}`$ pN. Merely by continuing the formulae developed for negative, i.e. resistive values of $`F_x`$, analytically through $`F_x=\mathrm{\hspace{0.17em}0}`$ to positive values, FK obtained an apparently reasonable description of the observed ‘acceleration ratios,’ $`V(F_x=\mathrm{\hspace{0.17em}5}\text{pN})/V(F_x=\mathrm{\hspace{0.17em}0})`$. These varied from about 3.0 to 1.4 as \[ATP\] increased from 5 $`\mu `$M to 1 mM. The more recent experiments, however, do not confirm these results lan:asb ; blo:asb ; car:cro ; on the contrary, for \[ATP\] $`50\mu `$M the data indicate no acceleration or even a deceleration, for $`F_x`$ up to $`10`$-15 pN, i.e., $`V(F_x>\mathrm{\hspace{0.17em}0})V(F_x=\mathrm{\hspace{0.17em}0})`$.
But, irrespective of the experimental observations, an examination of the geometry entailed in a standard single-bead assay with an assisting load, as illustrated in Figure 2 svo:blo ,
reveals that switching from a resistive to an assisting load should not be described merely by the change in sign of a scalar load force. Rather the true vectorial character of the force $`𝑭`$, acting at the point of attachment, P, of the tether to the body of the motor should be recognized. Even if sideways, $`F_y0`$, components of $`𝑭`$ may be neglected — although Block and coworkers lan:asb ; blo:asb have imposed transverse loads — a satisfactory description of the motor operation should seek to understand the mean velocity $`V`$, and likewise the dispersion $`D`$, as functions of $`F_x`$ and $`F_z`$.
Clearly, it would be advantageous experimentally to vary $`F_z`$ independently of $`F_x`$. With current optical trap technology this is not readily accomplished (since the traps are rather soft along the perpendicular or $`z`$-axis). To some degree, however, the issue can be addressed experimentally by employing beads of varying diameters. Thus, by reference to Figure 2, one sees that increasing the bead radius, $`R`$, increases $`F_z`$ relative to a controlled value of $`F_x`$ via
$$\frac{F_z}{F_x}=\mathrm{cot}\mathrm{\Theta }=\frac{R+\mathrm{\Delta }z}{[(l\mathrm{\Delta }z)(2R+l+\mathrm{\Delta }z)]^{1/2}}\sqrt{\frac{R}{2l}},$$
(2)
in which $`l`$ is the length of the tether while $`\mathrm{\Delta }z^2`$ represents the mean squared thermal fluctuation of the bead in the $`z`$ direction svo:blo .
A decade ago, however, Howard stressed the vectorial nature of the load and the significance of independently observing the perpendicular or vertical component, $`F_{}F_z`$, and measuring its effects on $`V`$. To this end Gittes, Meyhöfer, Baek and Howard (GMBH) git:mey devised an ingenious microtubule (MT) buckling experiment.
In their protocol the minus end of an MT is chemically attached to a glass substrate while the plus end of the MT can interact freely with a sparse field of kinesin motors tethered randomly to the same substrate. Then, as the free end of the MT diffusively wanders across the surface, a single kinesin molecule will occasionally encounter the MT and bind to it. In the presence of ATP it then proceeds to step along the MT towards the plus end: see Fig. 3. As the motor moves, it will eventually start to buckle the MT. By recording the buckling process and analyzing the shapes of the MT, the two components of the force, $`𝑭=(F_x,F_z)`$, in the plane of the substrate can be determined. Likewise the varying speed, $`V(t)`$, along the MT can be measured. The experiment is not easy and the resulting data are quite noisy: nevertheless, GMBH felt able to conclude that the vertical component, $`F_z`$, resulted in an increase in the overall mean speed.
In this article we extend the previous quantitative mechanochemical kinetic approach fis:kol ; fis:kol2 ; kol:fis in order to explicitly account for the vectorial character of the load force in analyzing the stochastic motions of motor proteins. The resulting formalism has been applied to the recent experiments of Block and coworkers on kinesin fis:kim ; kim:fis . Contrary to their model description, which implies sideways (or transverse) motions of the motor while stepping, we find no cause to invoke displacements of the point of attachment outside the $`(x,z)`$ plane. However, our analysis indicates that a kinesin motor ‘crouches’ on binding ATP, that is, the point of attachment of the tether moves downwards towards the microtubule track, apparently by about 0.5 to 0.8 nm before rapidly completing a unitary ‘swing’ or diffusive step of close to 8.0 nm transferring the motor to the next binding site on the track car:cro ; nis:mut .
Using the corresponding fitted kinetic parameters, we have revisited the Gittes et al. microtubule buckling experiment git:mey . As we explain below, our analysis avoids a simplifying assumption made in their discussion and seems qualitatively consistent with the buckling data; however, it indicates that $`V(F_x,F_z;\text{[ATP]})`$ generally decreases when $`F_z`$ increases (at fixed $`F_x`$ and \[ATP\]) rather than increases as GMBH argued. Nevertheless, one may note that in the buckling experiments the motor moves progressively away from a curved and, hence, stressed region of the microtubule. This raises the interesting question as to the degree to which the curvature of a microtubule protofilament might affect the motility of kinesin.
## 2 Intermediate States and Substeps
The modeling of motor protein motility may be carried out at different levels how ; lei:hus ; jul:adj ; kol:wid ; bus:kel . For concreteness and relative simplicity, we will focus on motors powered by the hydrolysis of ATP and moving like kinesin via steps of fixed size. Then there are specific binding sites located at positions $`x=ld`$ $`(l=0,\pm 1,\pm 2,\mathrm{})`$ along the periodic track where $`d`$ is the step distance corresponding to the track periodicity. For a microtubule one has $`d8.2`$ nm, representing the size of a tubulin dimer. In the absence of molecular fuel, $`\text{[ATP]}=0`$, the motor will be bound at a specific site $`l`$ in a nucleotide-free state to be labeled $`0_l`$. We will overlook the retention of ADP in the weakly bound head of kinesin how ; sch:cla and neglect the rates of spontaneous dissociation from the track: but see fis:kol2 ; kol:fis2 .
When, accepting the evidence for “tight coupling” how , the motor translocates from site $`l`$ to the adjacent site $`l+1`$, by processing a single fuel molecule, it undergoes, in the simplest linear reaction sequence considered here, a series of $`N`$ intermediate mechanochemical transitions from states $`j=0_l`$ to $`1_l`$ to $`\mathrm{}`$ to $`(N1)_l`$ to $`N_l0_{l+1}`$. When a motor is powered by the hydrolysis of ATP, the biochemical evidence indicates that $`N=4`$ distinct enzymatic states should be recognized how ; however, the degree to which these are linked to significantly different mechanical states is a significant object of investigation. Clearly, the simplest nontrivial model requires $`N=2`$ states: the transition from state $`0_l`$ to $`1_l`$ then corresponds to the binding of ATP which is followed by the hydrolysis process that is completed, with loss of ADP and P$`_\text{i}`$, when the motor moves on to state $`2_l0_{l+1}`$. Again for simplicity, we overlook the coupling between two distinct heads of a motor as entailed in the so-called hand-over-hand motion now well established for kinesin: see, e.g. sch:cla ; yil:tom ; klu:hoe:gil ; kol:phi . Nevertheless, this can be accommodated formally simply by doubling the nominal step size $`d`$ and allowing for twice as many intermediate states to describe periodic double-steps. (See, e.g., the analysis used in kol:fis for myosin V which, indeed, exhibits steps of fluctuating size.)
In the simplest, traditional chemical kinetic pictures fis:kol ; lei:hus one introduces transition rates, forwards and backwards, between states along a reaction pathway: thus we will associate rates $`u_j`$ and $`w_j`$ with the transitions from state $`j_l`$ forward to state $`(j`$$`+`$$`1)_l`$ and backward to state $`(j`$$``$$`1)_l`$, respectively. Note that owing to the periodicity of the stepping process, the rates $`u_j`$ and $`w_j`$ are independent of $`l`$.
However, to describe the action of motor proteins under variable loads it is essential to describe the dependence of the various rates on the loads imposed, or, complementarily, on the stresses induced: and how the influence of the load is distributed over the different mechanochemcial transitions should be a prime question fis:kol . Since purely chemical steps are typically very fast on the scale of the overall mechanical movements, a rather basic theoretical picture jul:adj postulates a distinct free energy surface, depending on molecular coordinates, for each separate mechanically fluctuating biochemical state. Motions occur by diffusion through the multidimensional configurational space while the probability of a formally instantaneous “chemical jump” from one surface to another varies with the configuration. One may argue, however kol:wid , that for practical purposes this scheme may be reduced to an effective mechanochemical description in which points along the traditional chemical “reaction coordinate” are, for a motor protein, simply identified with specific positions of the motor along the linear track.
By this route, as previously adopted fis:kol2 ; kol:fis , one may, indeed, hope to identify substeps, $`d_0`$, $`d_1`$, $`\mathrm{}`$, in which the motor moves from a binding site, say, at $`x_0`$ to the next site $`x_0+d`$ via intermediate locations $`x_1=x_0+d_0`$, $`x_2=x_0+d_0+d_1`$, $`\mathrm{}`$, with $`_{j=0}^{N1}d_j=d`$. Likewise, within the chemical picture, one may identify successive, unstable “transition states,” say $`j_l^+j_{l+1}^{}`$, along the reaction coordinate each lying between the locally stable intermediate states $`j_l`$ and $`j_{l+1}`$. In mapping the reaction coordinate onto the track, one may then decompose a substep $`d_j`$ from $`x_j`$ to $`x_{j+1}`$ into the sum $`(d_j^++d_{j+1}^{})`$ thereby locating the transition state at $`x_j^+=x_j+d_j^+=x_{j+1}d_{j+1}^{}x_{j+1}^{}`$: see Fig. 1 of fis:kol2 for depictions of such mappings with $`N=2`$ and $`N=4`$, as fitted to data for kinesin under resisting $`(F_x<0)`$ loads vis:sch ; sch:vis .
It is rather clear that in such a treatment, in which the reaction pathway is assumed to be parallel to the $`x`$-axis, only the longitudinal component, $`F_x`$, of the load will be coupled to displacements of the motor and so affect the rate constants. However, the initial substep predicted by this approach $`(d_0=1.8`$-2.1 nm$`)`$ appears to be inconsistent with high-resolution spatio-temporal observations of individual and averaged steps fis:kim ; kim:fis ; nis:mut ; car:cro . Furthermore, this approach has failed to account satisfactorily for velocity measurements of kinesin under reverse loading lan:asb ; blo:asb ; car:cro . Accordingly, apart from general theoretical considerations, it seems important to extend the previous discussions.
Evidently, a significant next step in considering the load-dependence is to recognize the vectorial character of the imposed force by allowing states along the reaction pathway to be mapped onto movements of the motor within the full three-dimensional space of the track, motor, tether, and cargo: see Fig. 1. More concretely we may hope to follow the course of the point of attachment P $`𝒓(t)=[x(t),y(t),z(t)]`$ of the tether to the motor as steps are taken: see Figs. 1 and 2. As varying values, $`z(t)`$ and $`y(t)`$, arise one can view the motor as moving up or down or swinging from side to side when it walks along the track from a binding site $`l`$ to the next binding site $`l+1`$. The simplest ($`N=\mathrm{\hspace{0.17em}2}`$)-state model embodying this concept is illustrated in Fig. 4.
It has been supposed, first, that the motion of P may be regarded as confined to the $`(x,z)`$ plane, i.e., that sideways (or transverse) motions may be neglected. Of course this should be subject to experimental test: but for kinesin lan:asb ; blo:asb ; kim:fis it proves quite adequate kim:fis .
Then in Fig. 4 the rectangular boxes labeled “0” and “$`20`$” represent bound, nucleotide-free states at successive sites $`l`$ and $`l+1`$, along the track. The circle labeled “1” denotes the ATP bound state; this state may also include subsequent ADP-P$`_\text{i}`$ and ADP biochemical states. This $`N=\mathrm{\hspace{0.17em}2}`$ level model, however, cannot distinguish situations in which the hydrolysis, ADP, and P$`_\text{i}`$-release chemical reactions take place physically at (or close to) the mechanical state “1” or close to “2”. (And recall also the remark above concerning hand-over-hand motion.)
The (unlabeled) crosses in Fig. 4 represent the location of the two transition states, $`0^+1^{}`$ and $`1^+2^{}`$. Also introduced in the figure are the dimensionless load distribution factors $`\theta _j^{x+}`$, $`\theta _j^x`$, $`\theta _j^{z+}`$ and $`\theta _j^z`$ (following fis:kol ; fis:kol2 ). These turn out to be crucial fitting parameters in describing experimental data. Evidently they serve to specify the substeps via
$$d_j^\pm =\theta _j^{x\pm }d\text{and}d_j=(\theta _j^{x+}+\theta _{j+1}^x)d,$$
(3)
while vertical (or perpendicular) displacements of the point of attachment are specified by
$$\mathrm{\Delta }z_j^\pm =\theta _j^{z\pm }d\text{and}\mathrm{\Delta }z_j=(\theta _j^{z+}+\theta _{j+1}^z)d,$$
(4)
with, of course, similar expression for $`\mathrm{\Delta }y^\pm `$ in terms of $`\theta _j^{y\pm }`$. Since the stepping is periodic one must finally have
$$\underset{j=0}{\overset{N1}{}}(\theta _j^{y+}+\theta _j^y)=\underset{j=0}{\overset{N1}{}}(\theta _j^{z+}+\theta _j^z)=0.$$
(5)
## 3 Load-dependence of Transition Rates
Now, under any external load, $`𝑭`$, the rate constants, $`u_j`$ and $`w_j`$, must change. But how? To answer, let us accpet a set of “free-energy landscapes”, for distinct biochemical states each depending on some set of molecular configurations jul:adj . Then one can contemplate a “reduced landscape” in which the set of distinct landscapes is, in effect, combined by identifying a traditional reaction coordinate, while at the same time retaining a particular mechanical coordinate, like the displacement, $`x`$, of the motor along its track bus:kel . (One may, incidently, note a critique fis of this general approach.) It is natural in the present case, as we have demonstrated, to extend the retained mechanical coordinate to be the vector $`𝒓`$ specifying, in real space, the point of attachment, P, of the tether to the motor. Then, if one supposes the reaction pathway can, at least on average, be mapped on to the motion of P, one may dispense with any explicit consideration of the reaction coordinate. Thus we are led to postulate a more-or-less well defined, presumably smooth and differentiable free-energy function, say $`\mathrm{\Phi }(𝒓)`$, which in the absence of a load (i.e., $`𝑭=0`$) respects the periodicity of the track so that
$$\mathrm{\Phi }(𝒓)=\mathrm{\Phi }(𝒓+d\widehat{𝒙}),$$
(6)
where $`\widehat{𝒙}`$ is a unit vector parallel to the $`x`$ axis. (See also how Chaps 15, 16, etc., and Sachs and Lecar sac:lec .)
The various mechanochemical states $`0_l`$, $`1_l`$, $`\mathrm{}`$, $`j_l`$, $`\mathrm{}`$ are then to be identified with corresponding valleys or potential wells, i.e., minima in $`\mathrm{\Phi }(𝒓)`$, located, say at $`𝒓_j+ld\widehat{𝒙}`$ for $`j=0,1,\mathrm{}`$ (mod $`N`$) and $`l=0,\pm 1,\pm 2,\mathrm{}`$. Following the traditional chemical picture, successive valleys along the reaction path will be linked via free energy barriers corresponding to the respective transition states. These, in turn, will be represented, by cols (or passes or saddles) in $`\mathrm{\Phi }(𝒓)`$ at points $`𝒓_j^+𝒓_{j+1}^{}`$.
When a force $`𝑭`$ is exerted on the motor’s tether the (now reduced) free energy landscape must be deformed. It is then most reasonable (but surely not fully ‘realistic’) to suppose that
$$\mathrm{\Phi }(𝒓;𝑭)=\mathrm{\Phi }(𝒓)𝑭\mathbf{}𝒓.$$
(7)
Now, once again in accord with traditional chemical reaction rate theories (see, e.g. how ), the rate $`u_j`$ of the forward reaction from a state $`j`$ to $`j+1`$ will be dominated by the Boltzmann factor for overcoming the corresponding barrier; and the same goes for the backwards reaction to $`j1`$. Thus, if $`\mathrm{\Phi }_j(𝑭)`$ is the free energy level at the bottom of the $`j`$th valley while $`\mathrm{\Phi }_j^+(𝑭)`$ and $`\mathrm{\Phi }_j^{}(𝑭)`$ are the height of the associated forward and rearward col, respectively, we may suppose
$`u_j(𝑭)`$ $``$ $`e^{[\mathrm{\Phi }_j^+(𝑭)\mathrm{\Phi }_j(𝑭)]/k_\text{B}T},`$
$`w_j(𝑭)`$ $``$ $`e^{[\mathrm{\Phi }_j^{}(𝑭)\mathrm{\Phi }_j(𝑭)]/k_\text{B}T}.`$ (8)
To proceed so as to obtain the force-dependence of the rates up to quadratic order in $`𝑭`$ on the basis of this landscape picture, consider first the forward rate constant $`u_j`$. In the absence of the load the corresponding valley is located at $`𝒓_j`$ while the appropriate col (describing the transition state $`j^+`$) is at $`𝒓_j^+𝒓_j+𝜽_j^+d`$ where we have introduced the forward load distribution vector $`𝜽_j^+=(\theta _j^{x+},\theta _j^{y+},\theta _j^{z+})`$. Expansion of the free energy in the valley then yields
$$\mathrm{\Phi }(𝒓)=\mathrm{\Phi }_j+\frac{1}{2}(𝒓𝒓_j)\mathbf{}\text{A}_j\mathbf{}(𝒓𝒓_j)+O(|𝒓𝒓_j|^3),$$
(9)
where $`\mathrm{\Phi }_j\mathrm{\Phi }(𝒓_j)`$ while A<sub>j</sub> is a positive definite symmetric $`3`$$`\times `$$`3`$ matrix with elements $`A_j^{\lambda \mu }(^2\mathrm{\Phi }/\lambda \mu )`$ evaluated at $`𝒓_j`$, where $`\lambda ,\mu =x,y,z`$. On the other hand, one may expand the free energy around the col in the form
$$\mathrm{\Phi }(𝒓)=\mathrm{\Phi }_j^++\frac{1}{2}(𝒓𝒓_j^+)\mathbf{}\text{A}_j^+\mathbf{}(𝒓𝒓_j^+)+\mathrm{},$$
(10)
where $`\mathrm{\Phi }_j^+\mathrm{\Phi }(𝒓_j^+)`$ and A$`{}_{}{}^{+}{}_{j}{}^{}`$ is again a $`3`$$`\times `$$`3`$ matrix with elements $`(^2\mathrm{\Phi }/\lambda \mu )`$ but now evaluated at $`𝒓_j^+`$. Note that by the character of a col or saddle point, the matrix A$`{}_{}{}^{+}{}_{j}{}^{}`$, which is identical to A$`{}_{}{}^{}{}_{j+1}{}^{}`$, must be indefinite with at least one negative eigenvalue. The corresponding eigenvector identifies the optimal direction of the reaction, projected into $`𝒓`$-space as the system enters and leaves the transition state. It will be appropriate for us to assume that the remaining eigenvalues are positive.
The rate constant $`u_j^0`$ for $`𝑭=0`$ is then proportional to $`\mathrm{exp}(\mathrm{\Delta }\mathrm{\Phi }_j^{+0}/k_\text{B}T)`$ where $`\mathrm{\Delta }\mathrm{\Phi }_j^{+0}=\mathrm{\Phi }(𝒓_j^+)\mathrm{\Phi }(𝒓_j)`$ is the barrier height. Under a vectorial load $`𝑭`$, however, the positions of both valley and col shift, the changes being proportional to $`𝑭`$ in leading order. The new positions may be found by solving the equations $`\mathrm{\Phi }(𝒓;𝑭)=0`$ using (9), (10) and (7). Thus one finds the free energy difference between the shifted valley and col to be
$`\mathrm{\Delta }\mathrm{\Phi }_j^+(𝑭)`$ $`=`$ $`\mathrm{\Phi }_j^+(𝑭)\mathrm{\Phi }_j(𝑭),`$
$`=`$ $`\mathrm{\Delta }\mathrm{\Phi }_j^{+0}𝑭\mathbf{}(𝒓_j^+𝒓_j)\frac{1}{2}𝑭\mathbf{}𝜼_j^+\mathbf{}𝑭+O(F^3),`$
in which appears the matrix
$$𝜼_j^+=(\text{A}_j^+)^1(\text{A}_j)^1,$$
(12)
where by our construction, the inverse matrices of A<sub>j</sub> and A$`{}_{}{}^{+}{}_{j}{}^{}`$ are well defined. Note, again, that $`𝒓_j^+𝒓_j=𝜽_j^+d`$.
By the same arguments, the free energy barrier for the reverse rate $`w_j`$ is determined by the transition state $`j^{}`$ located at $`𝒓_j^{}=𝒓_j𝜽_j^{}d`$ and is given by
$`\mathrm{\Delta }\mathrm{\Phi }_j^{}(𝑭)`$ $`=`$ $`\mathrm{\Phi }_j^{}(𝑭)\mathrm{\Phi }_j(𝑭),`$
$`=`$ $`\mathrm{\Delta }\mathrm{\Phi }_j^0+𝑭\mathbf{}𝜽_j^{}d+\frac{1}{2}𝑭\mathbf{}𝜼_j^{}\mathbf{}𝑭+O(F^3),`$
where $`𝜽_j^{}=(\theta _j^x,\theta _j^y,\theta _j^z)`$ is the reverse load distribution vector while
$$𝜼_j^{}=(\text{A}_j^{})^1(\text{A}_j)^1,$$
(14)
in which $`\text{A}_j^{}\text{A}_{j1}^+`$ is the matrix of the second derivatives of $`\mathrm{\Phi }`$ evaluated at $`𝒓_j^{}𝒓_{j1}^+`$.
Finally, we may express the load-dependence of the rate constants as
$`u_j(𝑭)`$ $`=`$ $`u_j^0\mathrm{exp}\{+[𝜽_j^+\mathbf{}𝑭d+\frac{1}{2}𝑭\mathbf{}𝜼_j^+\mathbf{}𝑭+O(F^3)]/k_\text{B}T\},`$
$`w_j(𝑭)`$ $`=`$ $`w_j^0\mathrm{exp}\{[𝜽_j^{}\mathbf{}𝑭d+\frac{1}{2}𝑭\mathbf{}𝜼_j^{}\mathbf{}𝑭+O(F^3)]/k_\text{B}T\}.`$
Evidently, the load distribution vectors $`𝜽_j^+`$ and $`𝜽_j^{}`$ simply generalize the previous load distribution scalars fis:kol ; fis:kol2 ; kol:fis and serve to locate the intermediate mechanochemical states and determine the vectorial character of the substeps linking them. By the underlying periodicity they must satisfy
$$\underset{j=0}{\overset{N1}{}}(𝜽_j^++𝜽_j^{})=\widehat{𝒙},$$
(17)
which simply summarizes (3) and (5).
Similarly, the matrices $`𝜼_j^+`$ and $`𝜼_j^{}`$ serve to measure the relative compliances of the various transition states. We may note that, by (12) and (14), the consecutive differences
$$𝜼_j^+𝜼_{j+1}^{}=(\text{A}_{j+1})^1(\text{A}_j)^1,$$
(18)
are independent of the properties of the transition states; then, via periodicity, one concludes that
$$\underset{j=0}{\overset{N1}{}}(𝜼_j^+𝜼_j^{})=0.$$
(19)
In summary, one might be tempted to regard the load-dependence expressions (LABEL:eq.15) and (LABEL:eq.16) merely as phenomenological expansions in powers of $`𝑭d/k_\text{B}T`$; however, the crucial feature lies in the physical interpretation of the load distribution vectors $`𝜽_j^\pm `$ and the compliance matrices $`𝜼_j^\pm `$ which, in turn, yields the constraints embodied in (17) and (19).
## 4 Motility of Kinesin under Vectorial Loading
Our primary aim now, as an application of the foregoing analysis, is to reconsider the GMBH buckling experiment git:mey . As explained in the Introduction with the aid of Fig. 3, the crucial feature of that study was the direct determination of the perpendicular force $`F_{}F_z`$ and, independently, the longitudinal or parallel component, given by $`F_{}F_x`$ (since the load is always resistive from the perspective of the track). We will suppose, in the absence of contrary evidence, that the microtubule (MT) does not twist and that the kinesin, when it attaches to the MT, binds on to the nearest protofilament which, again we suppose for simplicity, does not spiral around the MT how . Then any externally generated transverse force components may be neglected: i.e., $`F_y=0`$.
In order to analyze the buckling data theoretically following GMBH, one needs an explicit expression for $`V(F_x,F_z;\text{[ATP]})`$, the velocity of the motor along the track as a function of $`F_x`$, $`F_z`$ at fixed \[ATP\]. In default of a better description, GMBH postulated a simple linear dependence of $`V`$ on $`F_x`$ and $`F_z`$. However, we may hope to do better by using the recent data of Block and coworkers lan:asb ; blo:asb (BASL) who imposed assisting and resisting (as well as sideways) loads. The BASL experiments studied squid kinesin whereas GMBH employed kinesin from bovine brains: undoubtedly this and other disparate experimental details should play some role. Nevertheless, experience suggests that the quantitative effects should not be great.
Accordingly, in Fig. 5
we present two-dimensional velocity-contour plots of conventional kinesin in the $`(F_x,F_z)`$ plane at saturating ATP concentration by using the simple $`(N`$$`=`$$`2)`$-state kinetic model as fitted to the BASL data. Details of our analysis will be discussed elsewhere kim:fis ; fis:kim2 . For completeness, however, we explain briefly here how these contour plots have been derived.
In terms of the general $`N=\mathrm{\hspace{0.17em}2}`$ expression fis:kol
$$V=d(u_0u_1w_0w_1)/(u_0+u_1+w_0+w_1),$$
(20)
the dependence of the velocity as a function of \[ATP\] under (effectively) zero load may be accounted for, following fis:kol2 , by
$$u_0^0=k_0^0\text{[ATP]},w_0^0=k_0^{}\text{[ATP]}/(1+\text{[ATP]}/c_0)^{1/2},$$
(21)
in which the final, square-root factor (introduced to take account of an ATP regeneration system and a related trend in the stall force) plays only a small role. Then, with $`d=8.2`$ nm a good description of the kinesin data of BASL is given by the rates
$`k_0^0`$ $`=`$ $`1.35\mu \text{M}^1\text{s}^1,w_1^0=5.0\text{s}^1,u_1^0=100\text{s}^1,`$
$`k_0^{}`$ $`=`$ $`2.04\times 10^3\mu \text{M}^1\text{s}^1,c_0=20\mu \text{M},`$ (22)
which, apart for the larger value of $`k_0^{}`$, are quite comparable to the original fits fis:kol2 based on the data of Visscher and coworkers vis:sch ; sch:vis (restricted to $`F_x0`$).
Now Block et al. blo:asb measured the velocity $`V`$ only as a function of the parallel component $`F_x`$ (even though for both $`F_x<0`$ and $`F_x>0`$). However, in order to use (20) with the $`𝑭`$-dependence expressions (LABEL:eq.15) and (LABEL:eq.16), one also needs to know the perpendicular component $`F_z`$. The route to unraveling this puzzle depends, as indicated in the Introduction, on a consideration of the geometry of the bead-tether-motor-track configuration as illustrated in Fig. 2. However, further properties of the motor and the experimental set-up may need to be accounted for. Concretely, we desire a formula, say $`F_z=_z(F_x)`$, that correlates the imposed (but unmeasured) perpendicular component $`F_z`$ with the observed parallel component $`F_x`$ fis:kim .
The most basic hypothesis — termed Model 0 kim:fis ; fis:kim2 — is to suppose, following Fig. 2, that the tether angle $`\mathrm{\Theta }`$ simply switches sign when $`F_x`$ passes through zero: this is expressed by
$$\text{Model 0: }F_z=F_x\mathrm{cot}\mathrm{\Theta }(F_x)=c_{}|F_x|,$$
(23)
where $`c_{}=|\mathrm{cot}\mathrm{\Theta }|`$ follows from (2). Reasonable values for beads of diameter $`2R=0.50`$ $`\mu `$m, a fluctuating scale $`\mathrm{\Delta }z=5`$ nm svo:blo and a tether length $`l=60`$ nm yield $`\mathrm{\Theta }35^{}`$ and $`c_{}1.44`$.
In practice, the thermal fluctuations already recognized by allowing for $`\mathrm{\Delta }z`$ in Fig. 2, come into play strongly when $`|F_x|0.5`$ pN. Accordingly, a more realistic hypothesis (which is tested further in Sec. 6 below) is embodied in the smoothed form
$$\text{Model I: }F_z=_z(F_x)=c_{}\sqrt{F_x^2+F_0^2},$$
(24)
in which the fluctuation amplitude $`F_0=0.3`$ pN proves appropriate in light of observed force fluctuations lan:asb ; blo:asb . Study of the BASL data, however, reveals an unexpected and strong asymmetry in the relation $`F_z=_z(F_x)`$ for kinesin kim:fis ; fis:kim2 . This can be represented effectively via an additive term as
$$\text{Model II:}F_z=c_{}\left[\sqrt{F_x^2+F_0^2}+\frac{F_1}{2}\left(1+\frac{F_x}{\sqrt{F_x^2+F_0^2}}\right)\right].$$
(25)
The amplitude $`F_1`$, which is found to be close to 2.0 pN, measures the asymmetry. Although the new term may appear as an artificial construct it can be interpreted rather directly in mechanical terms kim:fis ; fis:kim2 . Note, incidentally, that the factor in large parentheses merely represents a smoothed Heaviside step function vanishing rapidly for $`F_x<0`$.
On this basis, successful fits to the BASL data \[including measurements of the randomness $`r(F_x;\text{[ATP]})=D/dV`$\] are achieved with the load distribution vectors
$`𝜽_0^+`$ $`=`$ $`(0.120,\mathrm{\hspace{0.33em}0},0.043),𝜽_1^+=(0.032,\mathrm{\hspace{0.33em}0},0.026),`$
$`𝜽_0^{}`$ $`=`$ $`(0.950,\mathrm{\hspace{0.33em}0},\mathrm{\hspace{0.33em}0.100}),𝜽_1^{}=(0.102,\mathrm{\hspace{0.33em}0},0.031).`$
Of course, these satisfy (17) so that only six independent fitting parameters are entailed. The compliance matrices, $`𝜼_0^+`$ and $`𝜼_0^{}`$, do not need to be invoked although one might reasonably presume that they are not totally negligible in reality.
Now we may note, using (3) and (4), that these load distribution vectors imply a very small forward step of only $`d_00.1`$ to 0.2 nm (allowing for the fitting uncertainties) on binding ATP but a much larger downwards displacement, namely $`\mathrm{\Delta }z_0=0.5`$ to $`0.8`$ nm. In this sense, then, the motor appears to “crouch” before it completes the main forward step of magnitude $`d_18.0`$ to $`8.1`$ nm. It should also be remarked fis:kol2 ; kol:fis , that the relatively large positively directed value of $`𝜽_0^{}`$ means that it is the reverse rate $`w_0`$ that changes most under a resisting load and thereby leads to the motor stalling $`(V=0)`$.
With the parameters presented above, the contours plotted in Fig. 5 follow from (20). We may remark that when the level of ATP is lowered the form of the contours changes little qualitatively but the scale is reduced in accord with (20)-(22) which imply close-to Michaelis-Menton variation of $`V(\text{[ATP]})`$ at zero load. It should also be noted that the direct sampling of the $`(F_x,F_z)`$ plane by the Block and coworkers experiments vis:sch ; sch:vis ; lan:asb ; blo:asb is essentially confined to the locus specified by the $`F_z=_z(F_x)`$ relation and so, by Model I, is roughly given by $`F_z=c_{}|F_x|0.3`$ pN. In principle, the behavior further from this locus could be somewhat different than portrayed. In particular, the fact that the fitting was achieved without explicitly employing any compliance coefficients is responsible for the steep rise of the velocity contours in Fig. 5 as the stall force, $`𝑭_\text{S}5.9\widehat{𝒙}`$ pN, is approached. Thus, in the absence of the compliance matrices and further $`𝑭`$-dependent corrections, the stall condition $`V=0`$ combined with the periodicity constraint (17) for the load distribution vectors implies independence of $`F_y`$ and $`F_z`$. Some evidence on the likely magnitudes of the compliance effects is available from analysis of the transverse loading experiments of Block et al. blo:asb ; fis:kim2 . Shifts in the contours, away from the $`_z(F_x)`$ locus, as large at 10 to $`25\%`$ might be realized.
## 5 Buckling of a Microtubule by a Motor
As explained above, in the experiments of GMBH git:mey the minus end of a microtubule was clamped to a flat glass substrate on which kinesin molecules were sparsely tethered. When an individual motor encountered and bound to the thermally flexing MT, it moved towards the plus end and soon buckled the MT: see Fig. 3. Then in time sequence (at 30 frames per second), the successive shapes of the MT were recorded. From the progress of the motor along the contour of the flexing MT the time varying velocity, $`V(t)`$, could be found by fitting the observed displacement curves. This was plotted (see git:mey ) against the corresponding force components, $`F_x`$ and $`F_z`$, derived from the shapes. Thereby, GMBH generated displays somewhat resembling those shown in Fig. 6,
although with significant fluctuations (and digitizing noise) in all three variables $`V`$, $`F_x`$ and $`F_z`$. In the cases reported, the parallel components, $`F_{}=F_x`$, started close to the stall force $`F_\text{S}6.5`$ pN and decreased more or less steadily with time after buckling, while the perpendicular components, $`F_{}=F_z`$, increased from close to zero. In parallel, the velocities started from low values, close to stall, but increased, sometimes very steeply with changing force, and eventually, for $`F_{}`$ and $`F_{}`$ in the range 2 $`5`$ pN, passed through maxima as high as 850 to 1100 nm/s. Since these velocities correspond or exceed the highest normally seen in standard single-bead assays under low loads $`|F_x|1`$ pN, GMBH concluded that the primary action of a perpendicular load was to increase $`V`$.
The $`(V,F_{x,z})`$ plots in Fig. 6, which roughly mimic the GMBH observations, were generated from the velocity contours shown in Fig. 5 by postulating the dashed and dot-dashed force trajectories labelled A and B, respectively. While our $`(V,F_{x,z})`$ plots resemble those of GMBH, they attain much lower maximal speeds. This is clearly a consequence of the facts (a) that, as normal, $`V(F_x=F_z=0)800`$ nm/s and (b) that for $`V200`$ nm/s the contours of $`V`$ in Fig. 5 slope upwards to the right so that, contrary to the conclusion of GMBH, the predominant effect of increasing $`F_z`$ is to reduce the velocity. It should also be remarked that the close to $`45^{}`$ slope of the velocity contours for $`F_x0`$ directly reflects the observations blo:asb ; car:cro that even large assisting loads do not result in any significant increase in $`V`$.
Before discussing possible reasons for the disagreement with the buckling experiments, we ask what form the corresponding trajectories should take in the $`(F_x,F_z)`$ plane. We will learn that the putative trajectories A and B in Fig. 5 are not very plausible.
Now, the persistence length of a microtubule is about 6 mm how . Thus in analyzing their buckling data git:mey , GMBH considered an MT as a stiff, uniform rod with a measured flexual rigidity $`EI`$ (where $`E`$ is the Young’s modulus and $`I`$ is the moment of inertia of the cross-section). The theory of bending a uniform rod is well established lan:lif ; but we will sketch it briefly in order to understand how the force trajectory in the buckling experiment may be derived. This then enables one to calculate the velocity of the kinesin molecule as it moves along the MT under the resulting induced load.
We will now assume that the MT is clamped at the origin of the $`(x,z)`$ plane, while the motor is located close to the $`x`$-axis, say at $`(L_0,0)`$, where $`L_0`$ is the initial distance along the MT to the binding site in the clamped position: see Fig. 3. Note that the coordinates $`x`$ and $`z`$ here differ from those introduced in Fig. 1 where the $`x`$-axis was taken parallel to the MT axis. Similarly, we will always denote the parallel and perpendicular components of the force relative to the MT as $`F_{}`$ and $`F_{}`$, respectively.
When a two-component force, $`𝑭=(F_{},F_{})`$, is exerted on the MT by the tethered kinesin at $`(L_0,0)`$, the shape of the buckled MT (considered as a uniform rod) satisfies the beam equation which can be written how
$$\frac{d^2\theta }{ds^2}=\beta ^2\mathrm{sin}[\theta (s)\phi _F]\text{with}\beta ^2=F/EI,$$
(27)
where $`s`$ is the arc length measured along the MT from the clamped origin $`(0,0)`$ while $`\theta (s)`$ is the tangential angle of the rod at the point $`[x(s),z(s)]`$ with respect to the $`x`$ axis; in addition, $`\phi _F`$ denotes the angle of the applied force $`𝑭`$ (with respect to the $`x`$ axis) at the location of the kinesin which defines the pivot point, while $`F=\sqrt{F_{}^2+F_{}^2}`$ is the magnitude of the buckling force.
Assuming that the clamping of the MT is tight and that no torque is exerted at the pivot point, the solution of this equation must satisfy the boundary conditions
$$\theta (s=\mathrm{\hspace{0.17em}0})=0(d\theta /ds)_{s=L}=0,$$
(28)
where $`L`$ is the total (time varying) arc length to the pivot point from the origin at the clamped position. Furthermore, the pivot point is fixed at the initial location of kinesin, which itself does not move spatially, so that one must have
$`x(s=L)`$ $`=`$ $`{\displaystyle _0^L}\mathrm{cos}\theta (s)𝑑s=L_0,`$
$`z(s=L)`$ $`=`$ $`{\displaystyle _0^L}\mathrm{sin}\theta (s)=0.`$ (29)
Now, if the angle $`\theta (s)`$ remains small — as it will when the MT starts to buckle — one may accept the approximations $`(dx/ds)=\mathrm{cos}\theta 1`$ and $`(dz/ds)=\mathrm{sin}\theta \theta `$ which lead to
$$\theta (dz/dx),(d^2\theta /ds^2)(d^3z/ds^3)(d^3z/dx^3).$$
(30)
Expanding (27) up to first order then yields the linear equation
$$\frac{d^3z}{dx^3}+\beta ^2\frac{dz}{dx}=\beta ^2\phi _F,$$
(31)
for the displacement $`z(x)`$ of the MT from the $`x`$ axis. This is subject to the boundary conditions $`z(0)=z^{}(0)=z^{\prime \prime }(0)=z(L_0)=0`$ which then give
$$z(x)=\phi _F\mathbf{(}xL_0+L_0\mathrm{cos}\beta x\beta ^1\mathrm{sin}\beta x\mathbf{)}.$$
(32)
But by (29) this form implies that the buckling force, given by $`F=\beta ^2EI`$, must satisfy the equation
$$\mathrm{tan}\beta L_0=\beta L_04.493.$$
(33)
Note that in this small-$`\theta `$ approximation the magnitude of the force remains constant at $`F^{}EI(4.493/L_0)^2`$ as the rod initially buckles. This in turn leads to a force trajectory in the $`(F_x,F_z)`$ plane that is just a semicircle of radius $`F^{}`$. Of course, the small-angle approximation fails when $`F_{}`$ increases.
Exact solutions to the nonlinear equation (27) can be expressed via an elliptic integral how ; git:mey as
$$\beta s=_{\varphi _0}^{\varphi (s)}\frac{d\varphi }{\sqrt{1k^2\mathrm{sin}^2\varphi }},$$
(34)
in which the modulus $`k`$ and elliptic angle $`\varphi (s)`$ are related by
$$k\mathrm{sin}\varphi (s)=\mathrm{sin}\frac{1}{2}[\theta (s)\phi _F],k^2=\mathrm{sin}^2\frac{1}{2}(\theta _L\phi _F),$$
(35)
while $`\theta _L\theta (s=L)`$. Note that the solution represented by (34) already satisfies the boundary conditions (28). However, the initial elliptic angle $`\varphi _0`$ and the modulus $`k`$ must be found so as to reproduce the correct pivot relations (29) while $`\beta \sqrt{F}`$ determines the scale of the arc length. The parallel and perpendicular force components for an MT buckled to a contour length $`L>L_0`$ are then given by
$`F_{}`$ $`=`$ $`F\mathrm{cos}(\theta _L\phi _F)=F(12k^2),`$ (36)
$`F_{}`$ $`=`$ $`F\mathrm{sin}(\theta _L\phi _F)=2Fk\sqrt{1k^2}.`$ (37)
Finally, for a given $`L>L_0`$, one may integrate (29) numerically using (34) and (35) to obtain the actual trajectory of MT buckling in the $`(F_x,F_z)`$ plane. The dotted curve in Fig. 5 represents such a trajectory when the initial parallel force is set equal to the stall force $`F_\text{S}5.9`$ pN. For small $`F_{}=F_z`$, the trajectory approaches the circle described by the small-angle approximation.
The $`(V,F_{x,z})`$ buckling plot following from the calculated force trajectory (dotted curve in Fig. 5) is presented in Fig. 7.
Evidently, this differs significantly from the results of GMBH and from the postulated forms shown in Fig. 6. Indeed, the key experimental finding, namely, that the velocity, although noisy, exhibits a maximum above 850 nm/s at intermediate force levels is not reproduced: on the contrary one finds $`V650`$ nm/s in this region. On the other hand, the theory presented by GMBH, which simplifies the functional form of the velocity to a linear expression in the force (see git:mey Fig. 9c) generates qualitatively similar behavior to Fig. 7. Thus $`F_{}`$ falls steadily while $`V`$ always increases and $`F_{}`$ passes through a maximum: see git:mey Fig. 10.
In calculating the force trajectory for buckling an MT, we assumed that the clamping is tight and that no torque is applied by the motor protein at the pivot point. These assumptions seem reasonable based on the fits to the experimental data presented in git:mey Fig. 5. Nevertheless, as one can detect in this figure, the experimental data, especially at the beginning of the buckling event, show small deviations from the fitted curve near the pivot point; indeed, the microtubule appears to be slightly bent at the pivot point. This suggests that the motor protein may actually exert a torque on the MT which might lead to a rather different force trajectory.
Other possible sources of uncertainty in the experiments and their interpretation were discussed critically by GMBH, including clamp looseness, misalignment of the motor and the axis of clamping ($`z=0`$ in Fig. 3), protofilament number variations and slight, preformed bends in the microtubules, etc. Nevertheless, it is difficult to understand how the large maximal velocities observed, exceeding 900 nm/s, could be significantly in error. If this conclusion stands — and, clearly, experiments in which $`F_z`$ and $`F_x`$ can be controlled more directly and with fewer uncertainties are to be desired — a satisfactory explanation remains to be found. It is certain that the interaction of the heads, i.e., the motor domains with the microtubule plays a crucial role in kinesin motility: indeed, a recent normal mode analysis of simple protein network models zhe:don highlights this feature as a distinction separating kinesin from standard myosin and the F1 ATPase motor. It seems possible, therefore, that the stressed state of the microtubule in the buckling experiments could be a prime cause of velocity enhancement.
## 6 Tether Linkage under Load
Underlying our analysis fis:kim ; kim:fis ; fis:kim2 of the Block et al. experiments imposing reverse and transverse loads on kinesin lan:asb ; blo:asb are the expressions embodied in Models I and II for the perpendicular force, $`F_z=_z(F_x)`$, induced via the tether and bead by the longitudinal (or parallel) load: see (24) and (25) in Sec. 4. It is clearly of interest to attempt to test these forms against experiments that examine the transmission of force in the bead-tether-motor-track assembly. To that end we may utilize data obtained by Svoboda and Block svo:blo who studied the elasticity of the bead-kinesin-microtubule linkage.
These authors measured the ratio of the velocity $`V_\text{b}`$ of a bead trapped by optical tweezers to the velocity $`V_\text{s}`$ of the stage to which the MT was attached, in the presence of the nonhydrolyzable ATP analog AMP-PNP. This produces a rigorlike association between the kinesin motor and the MT. The ratio may then be expressed in terms of elasticities as
$$V_\text{b}/V_\text{s}=K_{\text{mot}}/(K_{\text{mot}}+K_\text{b}),$$
(38)
where $`K_{\text{mot}}`$ and $`K_\text{b}`$ are the stiffness of the bead-motor-track linkage and of the bead in the optical trap, respectively. The ratio was measured as a function of the displacement $`x_\text{b}`$ of the bead from the center of the trap. The individual observations at a laser power of 62.5 mW are presented in Fig. 8.
Although the data are noisy owing to the Brownian motion of the bead and to linkage heterogeneity svo:blo , the mean ratio rises slowly as the displacement increases from the smallest observed value at $`x_\text{b}=50`$ nm: see the solid circles which have been calculated by averaging the data binned in intervals of size 10 nm. We may then ask how well our models fit these data.
To proceed, note that the longitudinal force, $`F_x`$, applied by the optical trap for a bead displacement $`x_\text{b}`$ satisfies
$$F_x=K_\text{b}x_\text{b}=K_{\text{mot}}\mathrm{\Delta }l_x,$$
(39)
where $`\mathrm{\Delta }l_x`$ is the change under tension of the projection of the kinesin tether on to the MT or $`x`$ axis. By assuming, as is reasonable how , that the kinesin tether is rigid, we have $`\mathrm{\Delta }l_x=l\mathrm{sin}|\mathrm{\Theta }|`$ where $`l`$ is the length of the tether while $`\mathrm{\Theta }`$, here negative, is the angle of the tether from the vertical axis: see Fig. 2. After some algebra, one finds
$$\frac{V_\text{b}}{V_\text{s}}=\frac{x_\text{b}}{x_\text{b}+l\mathrm{sin}|\mathrm{\Theta }|},\mathrm{sin}\mathrm{\Theta }=\frac{F_x}{\sqrt{F_x^2+F_z^2}}.$$
(40)
Finally, we may use (39) in combination with the models (24) and (25) for $`_z(F_x)`$, to predict the ratio as a function of $`x_\text{b}`$. For the data in Fig. 8 one has $`K_\text{b}=0.03`$ pN/nm svo:blo . If we take $`l=40`$ nm, the models I and II yield perfectly satisfactory (if not very informative) fits to the data as evident from the solid and dashed lines in the figure. Note the two values assumed for the force fluctuation $`F_0`$ in (24) and (25). The tether length fitted here is shorter than what would be expected for a free kinesin molecule on the basis of photomicrographs how , say, $`l60`$ nm. However, when a bead is chemically bound to a kinesin molecule it seems likely that some length of the tether may also be attached to the surface of the bead. Thus there are no grounds here for questioning the adequacy of the tether-based models for $`_z(F_x)`$.
## 7 Summary
We have extended the previous simple mechanochemical kinetic models for motor protein motion fis:kol ; fis:kol2 ; kol:fis to accommodate a three-dimensional free energy landscape for the point of attachment of the tether to the body of a motor that moves under a vector load $`𝑭`$. The load-dependence of the various forward and reverse rates describing the mechanochemical progress of the motor as it takes a single overall forward step of size $`d`$, are embodied, to first order in $`𝑭`$, in a set of load distribution vectors $`𝜽_j^+`$ and $`𝜽_j^{}`$. These, in turn, relate directly to the forward substeps, $`d_j(d)`$, taken by the motor along its track and also serve to localize the corresponding intermediate transition states. In quadratic order in the components $`(F_x,F_y,F_z)`$ of $`𝑭`$, a set of compliance matrices, $`𝜼_j^+`$ and $`𝜼_j^{}`$, comes into play.
Even when perpendicular (or vertical) force components, $`F_z`$, are not purposefully imposed on a motor, as in standard single-bead optical-trap assays, consideration of the geometry of the bead-tether-motor-track configuration demonstrates that loads with $`F_z>0`$ are induced. Furthermore, in switching between resistive and assisting loads (parallel to the track) accounting for these vertical components proves imperative.
The general analysis has been used to reconsider the microtubule buckling experiments of Howard and coworkers git:mey in which the values of $`F_z`$ were measured. To that end the approach was first applied to the recent experiments of Block et al. blo:asb in which, in particular, assisting loads imposed on kinesin were found to have little if any accelerating effect. The resulting fits provide velocity contours in the $`(F_x,F_z)`$ plane for kinesin under specified \[ATP\]: these are needed to address the buckling data of Howard and collaborators. Incidentally, however, the analysis also indicates that, on initially binding ATP, a kinesin motor “crouches,” i.e. moves downwards closer to the microtubule, prior to hydrolyzing ATP and stepping directly forward by $``$$`8`$ nm. Further and more detailed discussion of the Block et al. experiments on kinesin blo:asb is provided elsewhere kim:fis ; fis:kim2 ; but, as shown here, the associated modelling of the transmission of tension via the bead-tether-kinesin-microtubule linkage is consistent with previous measurement by Svoboda and Block svo:blo .
On the basis of the buckling experiments git:mey Howard and coworkers concluded that a perpendicular force $`F_z=3`$-5 pN, together with a longitudinal resisting force, $`|F_x|`$, of similar magnitude, results in kinesin velocities in excess of 900 nm/s. Our analysis does not substantiate this inference suggesting instead that velocities no larger than say 600 nm/s or so, should have been observed. A resolution of this challenging discrepancy might reside in the influence of microtubule stress or curvature on kinesin motility.
###### Acknowledgements.
We are grateful to Steven Block for providing us with his laboratory’s experimental data for kinesin. Interactions with him, Jonathon Howard, Matthew Lang and Anatoly Kolomeisky have been appreciated. The support of the National Science Foundation (through Grant No. CHE 03-01101) is gratefully acknowledged. |
warning/0506/hep-ex0506001.html | ar5iv | text | # Measurement of the 𝒕𝒕̄ Production Cross Section in 𝒑𝒑̄ Collisions at √𝒔=1.96𝐓𝐞𝐕 Using Lepton Plus Jets Events with Semileptonic B Decays to Muons
## I Introduction
Top quark pair production in the standard model proceeds via either quark-antiquark annihilation or gluon-gluon fusion. At the Fermilab Tevatron collider, with a center-of-mass energy of $`1.96\mathrm{TeV}`$, quark-antiquark annihilation is expected to dominate. For a top mass of $`175\mathrm{GeV}/c^2`$, the calculated cross section is $`6.7_{0.9}^{+0.7}`$pb theory , and is approximately 0.2 pb smaller for each $`1\mathrm{GeV}/c^2`$ increase in the value of the top mass over the range $`170\mathrm{GeV}/c^2<M_{\mathrm{top}}<190\mathrm{GeV}/c^2`$.
Measurements of the cross section for top quark pair production provide a test of QCD, as well as the standard model decay $`tWb`$. Non-standard model production mechanisms, such as the production and decay of a heavy resonance into $`t\overline{t}`$ pairs Xtt , could enhance the measured cross section. Non-standard model top quark decays, such as the decay into supersymmetric particles susy , could suppress the measured value, for which a $`tWb`$ branching fraction of nearly 100% is assumed.
In this paper we report a measurement of the $`t\overline{t}`$ production cross section in $`p\overline{p}`$ collisions at $`\sqrt{s}=1.96\mathrm{TeV}`$ with the CDF II detector at the Fermilab Tevatron. The standard model decay $`tWb`$ of the top quark results in a final state from $`t\overline{t}`$ production of two $`W`$ bosons and two bottom quarks. We select events consistent with a decay of one of the $`W`$ bosons to an electron or muon plus a neutrino, both of which have large momentum transverse to the beam direction ($`P_\mathrm{T}`$). We refer to these high $`P_\mathrm{T}`$ electrons or muons as the “primary lepton”. The neutrino is undetected and results in an imbalance in transverse momentum. The imbalance is labeled “missing $`E_\mathrm{T}`$” ($`E/_T`$) since it is reconstructed based on the flow of energy in the calorimeter MET . The other $`W`$ boson in the event decays to a pair of quarks. The two quarks from the $`W`$ boson and the two $`b`$ quarks from the top decays hadronize and are observed as jets of charged and neutral particles. This mode is referred as $`W`$ plus jets. We take advantage of the semileptonic decay of $`b`$ hadrons to muons to identify final-state jets that result from hadronization of the bottom quarks. Such a technique, called “soft-lepton tagging” (SLT), is effective in reducing the background to the $`t\overline{t}`$ signal from $`W`$ boson produced in association with several hadronic jets with large transverse momentum. The production cross section is measured in events with three or more jets and at least one SLT tagged jet.
This measurement is complementary to other measurements from CDF II, which use secondary vertex tagging, kinematic fitting, or a combination of the two SVX kinPRD SVXkin . A forthcoming paper combXsec will present a combined cross section measurement based on the result of these four analyses.
Previous measurements RunI from Run I at the Tevatron have measured production cross sections statistically consistent with the standard model prediction. This and other Run II measurements are made at a slightly higher center of mass energy ($`1.96\mathrm{TeV}`$ vs. $`1.8\mathrm{TeV}`$) and with nearly twice as much integrated luminosity.
The organization of this paper is as follows: Section II reviews the detector systems relevant to this analysis. The trigger and event selection, the data and the Monte Carlo samples and the SLT tagging algorithm are described in Section III. The estimate of the background is presented in Section IV. The acceptance and the $`t\overline{t}`$ event tagging efficiency are described in Section V. The evaluation of the systematic uncertainties on the measurement is presented in Section VI. The $`t\overline{t}`$ production cross section measurement and the conclusions are presented in Section VII and Section VIII.
## II The CDF II Detector
The CDF II detector is described in detail in CDF , only the components relevant to this measurement are summarized here. The CDF II detector is a nearly azimuthally and forward-backward symmetric detector designed to study $`p\overline{p}`$ interactions at the Fermilab Tevatron. It consists of a magnetic spectrometer surrounded by calorimeters and muon chambers. An elevation view of the CDF II detector is shown in Figure 1.
Charged particles are tracked inside a 1.4 T solenoidal magnetic field by an 8-layer silicon strip detector covering radii from 1.5 cm to 28 cm, followed by the central outer tracker (COT), an open-cell drift chamber that provides up to 96 measurements of charged particle position over the radial region from 40 cm to 137 cm. The 96 COT measurements are arranged in 8 “superlayers” of 12 sense wires each alternating between axial and 2 stereo. The COT and the silicon detectors track charged particles for $`|\eta |<1`$ and $`|\eta |<2`$, respectively.
Surrounding the tracking system are electromagnetic and hadronic calorimeters, used to measure charged and neutral particle energies. The electromagnetic calorimeter is a lead-scintillator sandwich and the hadronic calorimeter is an iron-scintillator sandwich. Both calorimeters are segmented in azimuth and polar angle to provide directional information for the energy deposition. The segmentation varies with position on the detector and is 15 in azimuth by 0.1 units of $`\eta `$ in the central region ($`|\eta |<1.1`$). Segmentation in the plug region ($`1.1<|\eta |<3.6`$) is 7.5 up to $`|\eta |<2.1`$, and 15 for $`|\eta |>2.1`$ in azimuth, while ranging from 0.1 to 0.64 units of $`\eta `$ in pseudo-rapidity (a nearly constant 2.7 change in polar angle). The electromagnetic calorimeters are instrumented with proportional and scintillating strip detectors that measure the transverse profile of electromagnetic showers at a depth corresponding to the shower maximum.
Outside the central calorimeter are four layers of muon drift chambers covering $`|\eta |<0.6`$ (CMU). The calorimeter provides approximately 1 meter of steel shielding. Behind an additional 60 cm of steel in the central region sit an additional four layers of muon drift chambers (CMP) arranged in a box-shaped layout around the central detector. Central muon extension (CMX) chambers, which are arrayed in a conical geometry, provide muon detection for the region $`0.6<|\eta |<1`$ with four to eight layers of drift chambers, depending on polar angle. All the muon chambers measure the coordinates of hits in the drift direction, $`x`$, via a drift time measurement and a calibrated drift velocity. The CMU and the CMX also measure the longitudinal coordinate, $`z`$. The longitudinal coordinate is measured in the CMU by comparing the height of pulses, encoded in time-over-threshold, at opposite ends of the sense wire. In the CMX, the conical geometry provides a small stereo angle from which the $`z`$ coordinate of track segments can be determined. Reconstructed track segments have a minimum of three hits, and a maximum of four hits in the CMU and the CMP, and 8 hits in the CMX.
## III Data Sample and Event Selection
In this section we describe the collision data and the Monte Carlo samples used in this analysis. Section III.1 outlines the trigger system used for the initial selection of events from the $`p\overline{p}`$ collisions. Section III.2 describes the Monte Carlo samples used for acceptance and background studies. The selection of the $`W`$+jets datasets from the triggered data samples is presented in Section III.3. The $`t\overline{t}`$ signal is extracted from the $`W`$+jets events through the identification of candidate $`b`$ hadron semileptonic decays to muons. The algorithm for identifying these decays is summarized in Section III.4, and its application to the $`W`$+jets dataset is described in Section III.5.
This analysis is based on an integrated luminosity of $`194\pm 11`$ pb<sup>-1</sup> klimenko collected with the CDF II detector between March 2002 and August 2003 (175 pb<sup>-1</sup> with the CMX detector operational).
### III.1 $`𝒑\overline{𝒑}`$ Collision Data
CDF II employs a three-level trigger system, the first two consisting of special purpose hardware and the third consisting of a farm of commodity computers. Triggers for this analysis are based on selecting high transverse momentum electrons and muons. The electron sample is triggered as follows: At the first trigger level events are selected by requiring a track with $`P_\mathrm{T}>8\mathrm{GeV}/c`$ matched to an electromagnetic calorimeter tower with $`E_\mathrm{T}>8\mathrm{GeV}`$ and little energy in the hadronic calorimeter behind it. At the second trigger level, calorimeter energy clusters are assembled and the track found at the first level is matched to an electromagnetic cluster with $`E_\mathrm{T}>16\mathrm{GeV}`$. At the third level, offline reconstruction is performed and an electron candidate with $`E_\mathrm{T}>18\mathrm{GeV}`$ is required. The efficiency of the electron trigger is measured from $`Zee`$ data and found to be $`(96.2\pm 0.6)\%`$ WZPRD . The selection of the muon sample begins at the first trigger level with a track with $`P_\mathrm{T}>4\mathrm{GeV}/c`$ matched to hits in the CMU and the CMP chambers or a track with $`P_\mathrm{T}>8\mathrm{GeV}/c`$ matched to hits in the CMX chambers. At the second level, a track with $`P_\mathrm{T}>8\mathrm{GeV}/c`$ is required in the event for about 70% of the integrated luminosity, while for the remainder, triggers at the first level are fed directly to the third level. At the third trigger level, a reconstructed track with $`P_\mathrm{T}>18\mathrm{GeV}/c`$ is required to be matched to the muon chamber hits. The efficiency of the muon trigger, measured from $`Z\mu \mu `$ data, is $`(88.7\pm 0.7)\%`$ for CMU/CMP muons and $`(95.4\pm 0.4)\%`$ for CMX muons WZPRD .
### III.2 Monte Carlo Datasets
The detector acceptance of $`t\overline{t}`$ events is modeled using PYTHIA v6.2 Pythia and HERWIG v6.4 Herwig . These are leading-order event generators with parton showering to simulate radiation and fragmentation effects. The generators are used with the CTEQ5L parton distribution functions CTEQ5L . Decays of $`b`$ and $`c`$ hadrons are modeled using QQ v9.1 QQ . Estimates of backgrounds from diboson production ($`WW`$, $`WZ`$, $`ZZ`$) are derived using the ALPGEN generator Alpgen with parton showering provided by HERWIG. The background from single top production (eg. $`W^{}t\overline{b}`$) is simulated using PYTHIA. Samples of the remaining backgrounds are derived directly from the data as described in Section IV.
The detector simulation reproduces the response of the detector to particles produced in $`p\overline{p}`$ collisions. The same detector geometry database is used in both the simulation and the reconstruction, and tracking of particles through matter is performed with GEANT3 geant . The drift model for the COT uses a parametrization of a GARFIELD simulation garfield with parameters tuned to match COT collider data. The calorimeter simulation uses the GFLASH gflash parametrization package interfaced with GEANT3. The GFLASH parameters are tuned to test beam data for electrons and high-$`P_\mathrm{T}`$ pions and checked by comparing the calorimeter energy of isolated tracks in the collision data to their momenta as measured in the COT. Further details of the CDF II simulation can be found in sim .
### III.3 $`𝑾`$+Jets Dataset
From the inclusive lepton dataset produced by the electron and muon triggers described in Section III.1, we select events with an isolated electron $`E_\mathrm{T}`$ (muon $`P_\mathrm{T}`$) greater than $`20\mathrm{GeV}`$ and $`E/_T>20\mathrm{GeV}`$. The isolation $`I`$ of the electron or muon is defined as the calorimeter transverse energy in a cone of $`\mathrm{\Delta }R\sqrt{\mathrm{\Delta }\eta ^2+\mathrm{\Delta }\varphi ^2}<0.4`$ around the lepton (not including the lepton energy itself) divided by the $`E_\mathrm{T}`$ ($`P_\mathrm{T}`$) of the lepton. We require $`I<0.1`$. The $`W`$+jets dataset is categorized according to the number of jets with $`E_\mathrm{T}>15\mathrm{GeV}`$ and $`|\eta |<2.0`$. The decay of $`t\overline{t}`$ pairs gives rise to events with typically at least three such jets, while the $`W`$ plus one or two jet samples provide a control dataset with little signal contamination. Jets are identified using a fixed-cone algorithm with a cone size of 0.4 and are constrained to originate from the $`p\overline{p}`$ collision vertex. Their energies are corrected to account for detector response variations in $`\eta `$, calorimeter gain instability, and multiple interactions in an event. A complete description of $`W`$+jets event selection is given in kinPRD .
The $`W`$+jets dataset consists mainly of events of direct production of $`W`$ bosons with multiple jets. This amounts also to the largest background to $`t\overline{t}`$ signal. As a first stage of background reduction, we define a total event energy, $`H_\mathrm{T}`$, as the scalar sum of the electron $`E_\mathrm{T}`$ or muon $`P_\mathrm{T}`$, the event $`E/_T`$ and jet $`E_\mathrm{T}`$ for jets with $`E_\mathrm{T}>8\mathrm{GeV}`$ and $`|\eta |<2.4`$. Due to the large mass of the top quark, a $`t\overline{t}`$ event is expected to have a large $`H_\mathrm{T}`$ compared to a $`W`$ plus three or more jets event, as illustrated in Figure 2. We studied the expected amount of signal (S) and background (B) as a function of $`H_\mathrm{T}`$ using the PYTHIA Monte Carlo program to model the signal $`H_\mathrm{T}`$ distribution. Data is used to model the background $`H_\mathrm{T}`$ distribution. We optimized the selection of events by imposing a minimum $`H_\mathrm{T}`$ requirement which maximizes $`S/\sqrt{S+B}`$. We select events with $`H_\mathrm{T}>200\mathrm{GeV}`$, rejecting approximately 40% of the background while retaining more than 95% of the $`t\overline{t}`$ signal.
There are 337 $`W`$ plus three or more jet events with $`H_\mathrm{T}>200\mathrm{GeV}`$ in 194 pb<sup>-1</sup> of data, 115 from $`W\mu \nu `$ candidates and 222 from $`We\nu `$ candidates.
Even after the $`H_\mathrm{T}`$ cut, the expected $`S:B`$ in $`W`$ plus three or more jet events is only of order 1:3. To further improve the signal to background ratio, we identify events with one or more $`b`$-jets by searching inside jets for semileptonic decays of $`b`$ hadrons into muons.
### III.4 Soft Lepton Tagging Algorithm
Muon identification at CDF proceeds by extrapolating tracks found in the central tracker, through the calorimeter to the muon chambers, and matching them to track segments reconstructed in the muon chambers. Matching is done in the following observables: extrapolated position along the muon chamber drift direction ($`x`$), the longitudinal coordinate along the chamber wires ($`z`$) when such information is available, and the extrapolated slope compared to the slope of the reconstructed muon chamber track segment ($`\varphi _L`$). Tracks are paired with muon chamber track segments based on the best match in $`x`$ for those track segments that are within 50 cm of an extrapolated COT track. In what follows we refer to the difference between the extrapolated and measured positions in $`x`$ and $`z`$ as d$`x`$ and d$`z`$, respectively, and the extrapolated and measured slope as d$`\varphi _L`$. The distributions of these variables over an ensemble of events are referred to as the matching distributions. In addition to selection based on d$`x`$ and d$`z`$, the standard muon identification also requires consistency with minimum ionizing energy deposition in the calorimeters. However, in order to retain sensitivity for muons embedded in jets, the muon SLT algorithm makes full usage of the muon matching information without any requirement on the calorimeter information. The algorithm starts with high-quality reconstructed tracks in the COT, selected by requiring at least 24 axial and 24 stereo hits on the track. Some rejection for pion and kaon decays in flight is achieved by requiring that the impact parameter of the reconstructed track be less than 3 mm with respect to the beamline. The track is also required to originate within 60 cm in $`z`$ of the center of the detector. Only tracks passing these cuts and extrapolating within $`3\sigma (P_\mathrm{T})`$ in $`x`$ outside of the muon chambers, where $`\sigma (P_\mathrm{T})`$ is the multiple scattering width, are considered as muon candidates. Also, when a track extrapolates to greater than $`3\sigma (P_\mathrm{T})`$ in $`x`$ inside the muon chambers, but no muon chamber track segment is found, the track is rejected and not allowed to be paired to other muon chamber track segments.
Candidate muons are selected with the SLT algorithm by constructing a quantity $`L`$, based on a comparison of the measured matching variables with their expectations. To construct $`L`$ we first form a sum, $`Q`$, of individual $`\chi ^2`$ variables
$$Q=\underset{i=1}{\overset{n}{}}\frac{(X_i\mu _i)^2}{\sigma _i^2},$$
(1)
where $`\mu _i`$ and $`\sigma _i`$ are the expected mean and width of the distribution of matching variable $`X_i`$. The sum is taken over $`n`$ selected variables, as described below. $`L`$ is then constructed by normalizing $`Q`$ according to
$$L=\frac{(Qn)}{\sqrt{\mathrm{var}(Q)}},$$
(2)
where the variance var$`(Q)`$ is calculated using the full covariance matrix for the selected variables. The normalization is chosen to make $`L`$ independent of the number of variables $`n`$; note that the distribution of $`L`$ tends to a Gaussian centered at zero and with unitary width, for $`n`$ sufficiently large. The correlation coefficients between each pair of variables are measured from $`J/\psi \mu \mu `$ data. The calculation proceeds by comparing the variance of the sum with the sum of the variances of each pair of $`\chi ^2`$ variables in Equation 1. Since the values of the matching variables are either positive or negative, according to the local coordinate system, separate correlation coefficients are used for pairs that have same-sign and opposite-sign values.
The selected variables are the full set of matching variables, $`x,z,\varphi _L`$ in the CMU, CMP and CMX with the following two exceptions: The CMP chambers do not provide a measurement of the longitudinal coordinate $`z`$, and matching in $`\varphi _L`$ is not included for track segments in the muon chambers that have only three hits. Because of their significantly poorer resolution, track segments reconstructed in the CMU chambers with three hits are not used. Note that a muon that traverses both the CMU and the CMP chambers yields two sets of matching measurements in $`x`$ and $`\varphi _L`$ and one $`z`$ matching measurement, and are referred as CMUP muons. All available matching variables are used in the calculation of $`L`$ for a given muon candidate. By placing an appropriate cut on $`L`$, background is preferentially rejected because hadrons have broader matching distributions than muons since the track segments in the muon chambers from hadrons are generally a result of leakage of the hadronic shower.
The widths of the matching distributions that enter into $`L`$ are a combination of intrinsic resolution of the muon chambers and multiple scattering. The multiple scattering term varies inversely with $`P_\mathrm{T}`$ and is dominant at low $`P_\mathrm{T}`$. The expected widths of the matching distributions are based on measurements of muons from $`J/\psi `$ decays at low $`P_\mathrm{T}`$ (see Figure 3) and $`W`$ and $`Z`$ boson decays at high $`P_\mathrm{T}`$.
The mean values ($`\mu _i`$ in Equation 1) are typically zero, except for a small offset in the CMU d$`z`$. We parameterize the widths as a function of up to three variables: $`P_\mathrm{T},\eta `$ and $`\varphi `$. These variables describe to first order the effects of multiple scattering in the detector. For the CMU detector, $`P_\mathrm{T}`$ is sufficient since the material traversed by a muon candidate is approximately homogenous in $`\eta `$ and $`\varphi `$. The widths are parameterized with a second-order polynomial in $`1/P_\mathrm{T}`$ with an exponential term that describes the $`P_\mathrm{T}`$ range below $`3\mathrm{GeV}/c`$. For the CMP detector we parameterize the widths as functions of $`P_\mathrm{T}`$ and $`\varphi `$ to take into account the rectangular shape of the absorber outside the central calorimeter. For the CMX detector we use $`P_\mathrm{T}`$, $`\eta `$ and $`\varphi `$ to account for a number of irregularities in the amount of absorber between $`\eta =0.6`$ and $`\eta =1.0`$. The measurement of the widths of the matching distributions as functions of $`P_\mathrm{T}`$, overlayed with their fits, are shown in Figure 4.
Figure 5 (left) shows an example of the distribution of $`L`$ from $`J/\psi `$ decays. The number of variables used varies from two to five. Figure 5 (right) shows the efficiency of the SLT algorithm as a function of $`L`$ from $`J/\psi `$ data. The efficiency plateaus at about 85% for $`|L|3.5`$.
### III.5 Event Tagging
In this analysis we seek to identify semileptonic decays of $`b`$ hadrons inside jets in $`t\overline{t}`$ events. The transverse momentum spectrum of these muons, covering a broad range from a few $`\mathrm{GeV}/c`$ to over $`40\mathrm{GeV}/c`$, is shown in Figure 6 from the PYTHIA Monte Carlo sample. Within the $`W`$+jets dataset defined in Section III.3, we isolate a subset of events with at least one “taggable” track. A taggable track is defined as any track, distinct from the primary lepton, passing the track quality requirements described in Section III.4, with $`P_\mathrm{T}>3\mathrm{GeV}/c`$, within $`\mathrm{\Delta }R<0.6`$ of a jet axis and pointing to the muon chambers to within a 3$`\sigma `$ multiple scattering window (the $`\sigma `$ of the multiple scattering window is defined as the $`\sigma _{\mathrm{d}x}`$ shown in Figure 3). The $`z`$-coordinate of the track at the origin must be within 5 cm of the reconstructed event vertex (the vertex reconstruction is described in detail in SVX ). Jets are considered “SLT tagged” if they contain a taggable track, which is also attached to a track segment in the muon chambers and the resulting muon candidate has $`|L|<3.5`$.
A potentially large background arises from $`J/\psi `$ decay and sequential, double semi-leptonic $`bcs`$ decay, resulting in one lepton from the $`b`$ decay and an oppositely charged lepton from the $`c`$ decay. Therefore, events are rejected if the primary lepton is of opposite charge to a SLT muon tag and the invariant mass of the pair is less than $`5\mathrm{GeV}/c^2`$. Similarly, events are also rejected if the primary lepton is a muon that together with an oppositely-charged SLT muon tag forms an invariant mass between 8 and $`11\mathrm{GeV}/c^2`$ or 70 and $`110\mathrm{GeV}/c^2`$, consistent with an $`\mathrm{{\rm Y}}`$ or a $`Z`$ particle, respectively. The sequential decay cut and the $`\mathrm{{\rm Y}}`$ and $`Z`$ removal reduce the $`t\overline{t}`$ acceptance by less than 1%.
Events passing all event selection cuts that have at least one taggable track are referred to as the ‘pretag’ sample. There are 319 pretag events with three or more jets, 211 in which the primary lepton is an electron and 108 in which it is a muon. Out of these events we find 20 events with a SLT tag, 15 in which the primary lepton is an electron and 5 in which it is a muon. This set of events is the $`t\overline{t}`$ candidate sample from which we measure the $`t\overline{t}`$ production cross section in Section VII.
## IV Backgrounds
In this section we describe the evaluation of background events in the $`t\overline{t}`$ candidates sample. The background contributions are mostly evaluated directly from the data. The dominant background in this analysis is from $`W`$ plus jets events where one jet produces an SLT tag. The estimate of this background relies on our ability to predict the number of such SLT tags starting from the pretag sample. The prediction is based on the probability for a given track in a jet to yield an SLT tag, and is measured in $`\gamma `$+jets events. We then evaluate the systematic uncertainty on the $`W`$+jets background estimate by testing the predictive power of the measured probabilities in a variety of data samples. The $`W`$+jets background evaluation and its systematic uncertainty is described in Section IV.1. After $`W`$+jets production, the next largest background is due to QCD multi-jet events. The evaluation of the QCD multi-jet contribution also relies on tagging probabilities measured in $`\gamma `$+jets events. However, we must account for the possible difference between the tagging probabilities for the QCD events that populate the $`t\overline{t}`$ candidate sample because the $`E/_T`$ often comes from a mismeasured jet and not from a neutrino. The evaluation of the QCD background is described in Section IV.2. An additional small source of background is due to Drell-Yan events and is estimated from the data and described in Section IV.3. The remaining background contributions are relatively small and are evaluated using Monte Carlo samples, as described in Section IV.4.
### IV.1 Backgrounds from $`𝑾`$+jets
$`W`$ plus jets events enter the signal sample either when one of the jets is a $`b`$-jet or a $`c`$-jet with a semileptonic decay to a muon, or a light quark jet is misidentified as containing a semileptonic decay (“mistagged”). We refer to these background events as $`W`$+heavy flavor and $`W`$+“fakes”, respectively. $`W`$+heavy flavor events include $`Wb\overline{b}`$, $`Wc\overline{c}`$ and $`Wc`$ production. One way of estimating these backgrounds would be to use a Monte Carlo program, such as ALPGEN to predict the $`W`$+heavy flavor component, and the data to predict the $`W`$+“fakes” (because the data provides a more reliable measure of mistags than the simulation). However, to avoid double-counting, this would require an estimate of mistags that is uncontaminated by tags from heavy flavor. Instead we have chosen to estimate both background components directly from the data, and we test the accuracy of the prediction as described below. We measure the combined $`W`$+heavy flavor and $`W`$+“fakes” background by constructing a “tag matrix” that parameterizes the probability that a taggable track with a given $`P_\mathrm{T}`$, $`\eta `$ and $`\varphi `$, in a jet with $`E_\mathrm{T}>15\mathrm{GeV}`$, will satisfy the SLT tagging requirement described in Section III.5. The variables $`\eta `$ and $`\varphi `$ are measured at the outer radius of the COT with respect to the origin of the CDF II coordinate system. The tag matrix is constructed using jets in $`\gamma `$+jets events with one or more jets. The tag probability is approximately 0.7% per taggable track, and includes tags from both fakes and heavy flavor. The tag rate as a function of each of the matrix parameters (integrated over the remaining two) is shown in Figure 7. The features in the tag rate vs. $`\eta `$ and $`\varphi `$ plots are a result of calorimeter gaps and changes in the thickness of the absorber before the muon chambers. The matrix is binned to take account of these variations. The bottom right plot shows the tag rate as the function of the $`|L|`$ cut for each muon category. The tag rate is higher for the CMP-only muons due to the smaller amount of absorber material that results from cracks in the calorimeter where there is no coverage by the CMU chambers.
We apply the tag matrix to all pretag events in the signal region according to:
$$N_{\mathrm{predicted}}^{\mathrm{tag}}=\underset{\mathrm{events}}{}\left[1\underset{i=1}{\overset{N_{\mathrm{trk}}}{}}\left(1𝒫_i\right)\right],$$
(3)
where the sum runs over each event in the pretag sample, and the product is over each taggable track in the event. $`𝒫_i`$ is the probability from the tag matrix for tagging the $`i`$-track with parameters $`P_{T_i}`$, $`\eta _i`$ and $`\varphi _i`$. Note that the sum over the events in Equation 3 includes any $`t\overline{t}`$ events that are in the pretag sample. We correct for the resulting overestimate of the background at the final stage of the cross section calculation, since the correction depends on the measured tagging efficiency (see Section VII).
A fraction, $`F_{QCD}`$, of the events in the signal region are QCD events (such as $`b\overline{b}`$ or events in which a jet fakes an isolated lepton, see Section IV.2) for which the background is estimated separately. Therefore, we explicitly exclude their contribution to $`N_{\mathrm{predicted}}^{\mathrm{tag}}`$ and obtain the predicted number of tagged $`W`$+jets background events
$$N_{\mathrm{predicted}}^{Wj\mathrm{tag}}=(1F_{QCD})N_{\mathrm{predicted}}^{\mathrm{tag}}.$$
(4)
The estimated $`W`$+fakes and $`W`$+heavy flavor background is given in the third line of Table 1.
The above technique relies on the assumption that the tagging rate in jets of the $`\gamma `$+jets sample is a good model for the tagging rate of the jets in $`W`$+jets events. The assumption is plausible because the SLT tagging rate in generic jet events is largely due to fakes. We have studied the heavy flavor content of the tags in the $`\gamma `$+jets sample using the overlap sample between SLT tags and displaced vertex tags identified with the silicon tracker SVX . We find that (21.0$`\pm `$1.4)% of the tags in the $`\gamma `$+jets sample are from heavy flavor. We have used MADEVENT MadEvent to do generator-level comparisons of the heavy-flavor fractions of $`W`$+jets events with those from the $`\gamma `$ plus jets events that make up the tag matrix. We find that the $`\gamma `$+jets sample used to make the tag matrix has approximately 30% more heavy flavor than the $`W+3`$ jet events. Since SLT tags in $`\gamma `$+jets events are dominantly fakes, this difference affects the background prediction in $`W`$+jets events at only the few-percent level.
Given the limitations of a generator-level, matrix element Monte Carlo study of the heavy flavor content of the $`\gamma `$+jets and $`W`$+jets samples, we do not use the MADEVENT study to evaluate the systematic uncertainty on the background due to tagged $`W`$+jets events. Instead we test the accuracy of the tag matrix for predicting SLT muon tags by using it to predict the number of tags in a variety of samples with different heavy flavor content. We check $`Z`$ plus jets events, events triggered on a jet with $`E_\mathrm{T}`$ thresholds of 20, 50, 70 and 100 $`\mathrm{GeV}`$ (called Jet 20, Jet 50, Jet 70 and Jet 100), or triggered on four jets and the scalar sum of transverse energy in the detector (called SumET). We find that the matrix predicts the observed number of tags in each of these samples to within 10%, as shown in Figure 8, and we use this as the systematic uncertainty on the prediction from the tag matrix.
### IV.2 QCD Background
We refer to events with two or more jets in which the decay of a heavy-flavor hadron produces a high-$`P_\mathrm{T}`$ isolated lepton, or in which a jet fakes such a lepton, as QCD events. These events enter the sample when, in addition to the high-$`P_\mathrm{T}`$ isolated lepton, a muon from a heavy flavor decay gives an SLT tag, or there is a fake tag. We measure this background directly from the data.
To estimate the QCD component we first use the distribution of pretag events in the plane of $`E/_T`$ vs. isolation, $`I`$, of the primary lepton. We populate this plane with lepton plus jets events according to the event $`E/_T`$ and $`I`$. We consider four regions in the plane:
$`A:E/_T<15`$ $`I>0.2`$
$`B:E/_T<15`$ $`I<0.1`$
$`C:E/_T>20`$ $`I>0.2`$
$`D:E/_T>20`$ $`I<0.1`$
where Region D is the $`t\overline{t}`$ signal region. The distribution of events, with one or more jets, in the $`E/_T`$ vs. $`I`$ regions is shown in Figure 9.
In order to populate Regions A, B and C with only QCD events, we correct the number of events for the expected contamination of $`W`$+jets and $`t\overline{t}`$ events in those regions using expectations from PYTHIA $`t\overline{t}`$ and $`W`$ Monte Carlo simulations. The corrections range from less than 1% in electron plus one-jet events in Region A, to (57$`\pm `$15)% in Region B in muon plus three or more jet events.
Assuming that the variables $`E/_T`$ and $`I`$ are uncorrelated for the QCD background, the ratio of the number of QCD events in Region A to those in Region B should be the same as the ratio of the number of QCD events in Region C to those in Region D. Therefore we calculate the fraction of QCD events in Region D, $`F_{QCD}`$, as:
$$F_{QCD}=\frac{N_D^{QCD}}{N_D}|_{\mathrm{pretag}}=\frac{N_BN_C}{N_AN_D}|_{\mathrm{pretag}},$$
(5)
where $`N_D^{QCD}`$ is the total number of pretag QCD events in the signal region, and $`N_i`$ represent the number of events in region $`i`$. The measured fractions are shown in Table 2.
To estimate the number of tagged QCD events in the signal region, we multiply $`F_{QCD}`$ by the tagging probability for QCD events. However, this tagging probability is not necessarily given by the tag matrix probabilities which are designed for jets in $`W`$+jets events. Mismeasurement in the jet energies and differences in kinematics between $`W`$+jets and QCD events may affect the tagging probabilities. $`W`$+jets events have $`E/_T`$ from the undetected neutrino, whereas QCD events have $`E/_T`$ primarily from jet mismeasurement. Jet mismeasurement is correlated with fake tags due to energy leakage from the calorimeter through calorimeter gaps or incomplete absorption of the hadronic shower, both of which can result in track segments in the muon chambers. $`W`$+jets events have a primary lepton from the $`W`$ decay, whereas QCD events have a primary lepton that is either a fake or a result of a semileptonic decay of heavy flavor. The presence of a lepton from heavy flavor decay typically enhances the tag rate. Figure 10 shows the ratio of the number of measured tags in the Jet 20 sample to the number of tags predicted by the tag matrix as a function of $`E/_T`$. As expected, in QCD events with large $`E/_T`$ we find a tag rate significantly larger than that described by the tag matrix.
We find that the prediction of the tag matrix can be renormalized to properly account for the tag rates in QCD events with a single multiplicative factor, which we call $`k`$. We measure $`k`$ using events in region C by comparing the number of SLT tags found to the number predicted by the tag matrix. Since the signal region contains only isolated ($`I<0.1`$) primary leptons, we reject events in the measurement of $`k`$ in which the SLT tag is within $`\mathrm{\Delta }R<0.5`$ of the primary lepton. After this requirement we do not find any dependence of $`k`$ on the isolation of the primary lepton. Figure 11 shows the ratio of measured to predicted tags in events in region C as a function of $`H_\mathrm{T}`$. The tag rate above $`H_\mathrm{T}=200\mathrm{GeV}`$ is approximately flat and is not much different from the prediction of the tag matrix (dashed line in Figure 11). However, QCD events at lower $`H_\mathrm{T}`$ have a significantly different tag rate than that predicted by the tag matrix. As shown in Figure 12, $`F_{QCD}`$ also has an $`H_\mathrm{T}`$ dependence for events with 1 or 2 jets, but is flat within the statistical uncertainty for three or more jets.
The number of QCD background events is calculated as:
$$N_{QCD}=F_{QCD}kN_{\mathrm{predicted}}^{\mathrm{tag}},$$
(6)
where $`N_{\mathrm{predicted}}^{\mathrm{tag}}`$ is given in Equation 3 and the brackets represent the product of $`F_{QCD}`$ and $`k`$ convoluted with the $`H_\mathrm{T}`$ distribution of QCD events from region C. In the control region (1 and 2 jets), the fits of $`F_{QCD}`$ vs. $`H_\mathrm{T}`$, shown in Figure 12, are convoluted with $`k`$ vs. $`H_\mathrm{T}`$, shown in Figure 11. For events with three or more jets, since there is no visible $`H_\mathrm{T}`$ dependence for either $`F_{QCD}`$ or $`k`$, we simply multiply their average values for $`H_\mathrm{T}>200\mathrm{GeV}`$. Measured values of $`k`$ times $`F_{QCD}`$ are given in Table 2. The procedure by which $`F_{QCD}`$ is determined as a function of $`H_\mathrm{T}`$ is important because the ratios between the four different regions of the $`E/_T`$ and $`I`$ kinematic plane, calculated in separate ranges of $`H_\mathrm{T}`$ and then averaged, does not necessarily correspond to the same ratios taken while integrating over the full $`H_\mathrm{T}`$ range. The uncertainties on the 1 and 2 jet events are conservatively taken as the difference between the central value and the result of the straight product of $`F_{QCD}`$ and $`k`$. The straight product corresponds to ignoring the $`H_\mathrm{T}`$ dependence as well as any other variable’s dependence of $`F_{QCD}`$ and $`k`$.
### IV.3 Drell-Yan Background
Drell-Yan events can enter the sample when they are produced with jets and one muon is identified as the primary muon while the second muon is close enough to a jet to be tagged. Residual Drell-Yan background that is not removed by the dimuon and sequential decay rejection described in Section III.5, is estimated from the data. We use events inside the $`Z`$-mass window ($`76106\mathrm{GeV}/c^2`$), which are otherwise removed from the analysis by the $`Z`$-mass cut, to measure the number of events that would pass all our selection requirements including the SLT tag, $`N_{\mathrm{inside}}^{\mathrm{tags}}`$. Because of the limited sample size of $`Z`$+jets events, we use $`Z+0\mathrm{jet}`$ events without the $`E/_T`$ and $`H_\mathrm{T}`$ requirements to find the ratio of events outside the $`Z`$-mass window to those inside the window, $`R_{Z/\gamma ^{}}^{\mathrm{out}/\mathrm{in}}`$, and a first-order estimate of the number of expected Drell-Yan events outside of the $`Z`$-mass window is calculated as:
$$N_{DY}=N_{\mathrm{inside}}^{\mathrm{tags}}R_{Z/\gamma ^{}}^{\mathrm{out}/\mathrm{in}}.$$
(7)
This estimate assumes that $`R_{Z/\gamma ^{}}^{\mathrm{out}/\mathrm{in}}`$ does not depend on the number of jets in the event. We assign a systematic uncertainty of 33% for this assumption based on the largest deviation between $`R_{Z/\gamma ^{}}^{\mathrm{out}/\mathrm{in}}`$ for ALPGEN $`Z/\gamma ^{}`$ plus zero jet events compared with 1, 2 or $``$3 jets events.
The first-order estimate is then corrected by the relative efficiency inside and outside the $`Z`$-mass window of the $`E/_T`$, $`H_\mathrm{T}`$, and SLT-jet requirements, which we measure using $`Z/\gamma ^{}`$+jets Monte Carlo events. The Drell-Yan background estimates are listed in the sixth line of Table 1.
### IV.4 Other Backgrounds
Remaining background sources are due to $`WW`$, $`WZ`$, $`ZZ`$, $`Z\tau \tau `$ and single top production. Diboson events can enter the sample when there are two leptons from a $`Z`$ and/or a $`W`$ decay and jets. One lepton passes the primary lepton requirements while the second is available to pass the SLT requirement if it is close to a jet. The $`E/_T`$ in these events can either come from a $`W`$-boson decay or from an undetected lepton in a $`Z`$-boson decay. $`Z\tau \tau `$ events can enter the sample when the $`Z`$ is produced in association with jets and one $`\tau `$ decays to a high-$`P_\mathrm{T}`$ isolated electron or muon, while the second $`\tau `$ produces an SLT muon in its decay. Electroweak single top production gives rise to an event signature nearly identical to $`t\overline{t}`$ when there are additional jets from gluon radiation.
None of the above background sources are completely accounted for by the application of the tag matrix to the pretag event sample because these backgrounds have a significant source of muons from, for instance, $`W`$ and $`Z`$ decay. Therefore, we independently estimate their contributions to the background using Monte Carlo samples normalized to the cross sections referenced in Table 3. In modeling the SLT tagging of such events in the Monte Carlo samples, we explicitly exclude the mistag contribution which is taken in to account in the application of the tag matrix to the pretag sample. The background for each source is estimated as:
$$N_i=\sigma _iA_iϵ_{\mathrm{tag},i}\mathrm{dt},$$
(8)
where $`\sigma _i`$ is the theoretical cross section for the particular background source, $`A_i`$ is the acceptance for passing the pretag event selection, $`ϵ_{\mathrm{tag},i}`$ is the SLT tagging efficiency and $`\mathrm{dt}`$ is the integrated luminosity of the overall data sample. The expected background contributions are shown, as a function of jet multiplicity, in Table 3.
## V Total $`𝒕\overline{𝒕}`$ Acceptance
We factorize the efficiency for identifying $`t\overline{t}`$ events into the geometric times kinematic acceptance and the SLT tagging efficiency. The acceptance includes all the cuts described in Sections III.3 as well as the invariant mass cut described in Section III.5, and is evaluated assuming a top mass of 175 GeV/c<sup>2</sup>. The tagging efficiency is the efficiency for SLT-tagging at least one jet in events that pass the geometric and kinematic selection. We describe each piece below.
### V.1 Geometric and Kinematic Acceptance
The acceptance is measured in a combination of data and Monte Carlo simulations. Simulations are done using the PYTHIA Monte Carlo program Pythia . The primary lepton identification efficiency is measured in $`Z`$-boson decays acquired with a trigger that requires a single high-$`P_\mathrm{T}`$ electron or muon. The efficiency is measured using the lepton from the $`Z`$-boson decay that is unbiased by the trigger, and the identification efficiency in the Monte Carlo sample is scaled to that measured in the data kinPRD . The acceptance, as a function of the number of identified jets above $`15\mathrm{GeV}`$, is shown in Table 4. These numbers include the measured efficiencies of the high-$`P_\mathrm{T}`$ lepton triggers.
### V.2 SLT Efficiency
The efficiency for the reconstruction of the COT track is taken directly from Monte Carlo simulation. The reconstruction efficiencies of muon chamber track segments are also taken from the simulation and scaled to the values measured in the data using the lepton in $`Z`$-boson decays unbiased by the trigger. The muon identification efficiency of the SLT algorithm is measured in data using $`J/\psi `$ and $`Z`$ decays. We use events acquired with triggers that demand a single muon, and use only the muon not biased by the trigger. The efficiency is defined as the ratio of muons that satisfy the SLT tagging requirement over the number of taggable tracks attached to track segments in the muon chambers. In the calculation of the efficiency a background linear in invariant mass is subtracted from the $`J/\psi `$ and $`Z`$ peaks. The measured efficiency vs. $`P_\mathrm{T}`$ is shown in Figure 13 for muons with $`|\eta |<0.6`$ (CMU and/or CMP) and for muons with $`0.6|\eta |1.0`$ (CMX). The decrease in efficiency with increasing $`P_\mathrm{T}`$ is a result of non-Gaussian tails in the components of $`L`$.
Since the efficiency measurement is dominated by isolated muons, whereas the muons in $`b`$-jets tend to be surrounded by other tracks, we have studied the dependence of the efficiency on the number of tracks, $`N_{\mathrm{trk}}`$, above $`1\mathrm{GeV}/c`$ in a cone of $`\mathrm{\Delta }R=0.4`$ around the muon track. We find no significant efficiency loss, although the precision of the measurement is poor near $`N_{\mathrm{trk}}=6`$, the mean expected in $`t\overline{t}`$ events. We include a systematic uncertainty to account for this by fitting the efficiency vs. $`N_{\mathrm{trk}}`$ to a linear function and evaluating this function at the mean $`N_{\mathrm{trk}}`$ expected in $`t\overline{t}`$ events. The systematic uncertainty on the efficiency for at least one SLT tag in a $`t\overline{t}`$ event from this effect is +0%, $``$8%.
The detector simulation does not properly reproduce the non-Gaussian tails of the muon matching distributions. Therefore the measured efficiencies, shown in Figure 13, are applied directly to a generated muon in the Monte Carlo sample when evaluating the efficiency for tagging a $`t\overline{t}`$ event. This accounts for tagging of semileptonic heavy flavor decays in $`t\overline{t}`$ events (including charm decays from $`Wc\overline{s}`$). Events from $`t\overline{t}`$ can also be mistagged when a tag results from a fake muon or a decay-in-flight. We account for this effect in the tagging efficiency evaluation by allowing events that are not tagged by muons from heavy flavor decays to be tagged by other charged tracks using the tagging probabilities from the tag matrix, as described in Section IV. Since the heavy flavor component of the tagging efficiency has already been accounted for, the generic-track tagging probabilities are corrected downwards for the measured 21% heavy flavor component of the tag matrix (Section IV.1). The overall efficiency for finding one or more SLT tags in a $`t\overline{t}`$ event (“tagging efficiency”) is shown in Table 5. Mistags account for approximately 25% of the $`t\overline{t}`$ tagging efficiency. Because a small portion of the integrated luminosity was accumulated before the CMX was fully functional, we break the efficiency into pieces with and without the CMX. This is taken into account in the final cross section denominator. The total $`t\overline{t}`$ detection efficiency is the product of the acceptance and the tagging efficiency.
As noted above, the SLT efficiency has been parameterized using muons that tend to be isolated from other activity. To further check that this efficiency measurement is representative of muons in or near jets, we use a high-purity $`b\overline{b}`$ sample, derived from events triggered on $`8\mathrm{GeV}`$ electrons or muons. These events are enriched in semileptonic $`b`$-hadron decays. To select this sample, we require that the events have two jets above $`15\mathrm{GeV}`$. One jet must be within $`\mathrm{\Delta }R=0.4`$ of the primary electron or muon (the “lepton jet”). For jets associated with muons, the energy is corrected to account for the muon $`P_\mathrm{T}`$. The second jet (the “away jet”) in the event is chosen as the jet above $`15\mathrm{GeV}`$ with maximum separation in azimuth ($``$2 radians) from the lepton jet. Both jets are required to have a secondary vertex reconstructed and tagged by the SecVtx algorithm SVX (“SecVtx-tagged”). This results in a $`b\overline{b}`$ sample with a purity of approximately 95% Lannon . We measure the SLT acceptance times efficiency for semileptonic decays to muons in the away jet in a HERWIG dijet Monte Carlo sample. Monte Carlo events are subject to the same event selection, as described above, used for the $`b\overline{b}`$ data sample. The efficiency parametrization measured from the data is applied in the same way as in the $`t\overline{t}`$ Monte Carlo sample. The derived efficiency times acceptance per $`b`$-jet is applied to the data to predict the number of SLT tags in the away jet. There are 7726 SecVtx-tagged away jets in which the lepton jet is from a muon and 2233 in which it is from an electron. In these events we predict $`388\pm 54`$ tags in the away jet opposite a muon jet and $`116\pm 17`$ tags in the away-jet opposite an electron jet. We find 353 and 106 respectively. We conclude that the efficiency for SLT-tagging muons from semileptonic decays of heavy flavor in jets is well modeled by our simulation.
## VI Systematic Uncertainties
Systematic uncertainties in this analysis come from uncertainties in the Monte Carlo modeling of the acceptance, knowledge of the SLT tagging efficiency, the effect on the acceptance of the uncertainty on the jet energy calibration, uncertainties on the background predictions, and the uncertainty on the luminosity.
Uncertainties in the Monte Carlo modeling of acceptance include effects of parton distribution functions (PDFs), initial-state radiation (ISR), final-state radiation (FSR), and the calibration of the measured jet energy. These are estimated by comparing different choices for PDFs, varying ISR, FSR and the jet energy in the Monte Carlo programs and comparing the results from the PYTHIA generator with those from HERWIG. A complete description of the evaluation of these uncertainties appears in kinPRD . The total systematic uncertainty on the acceptance due to these factors is $`\pm `$6.1%. Possible variations of the lepton ID efficiency in events with multiple jets are an additional source of systematic uncertainty on the acceptance. We use a data to Monte Carlo scale factor for the lepton ID efficiency that is taken from $`Zee`$ and $`Z\mu \mu `$ data and Monte Carlo samples. These samples contain predominantly events with no jets. A 5% systematic uncertainty on the scale factor is estimated by convoluting the scale factor itself, measured as a function of $`\mathrm{\Delta }R`$ between the lepton and the nearest jet, with the $`\mathrm{\Delta }R`$ distribution of leptons in $``$3 jet $`t\overline{t}`$ events kinPRD . Adding the uncertainties in quadrature gives a total Monte Carlo modeling systematic uncertainty on the acceptance of $`\pm `$8.0% .
There are several factors that contribute to the systematic uncertainty on the SLT tagging efficiency. The uncertainty due to our knowledge of the $`P_\mathrm{T}`$ dependence is determined by varying the efficiency curves used in the $`t\overline{t}`$ Monte Carlo sample according to the upper and lower bands in Figures 13. We find that the tagging efficiency for $`t\overline{t}`$ changes by $`\pm `$1% from its central value. An additional source of systematic uncertainty for the tagging efficiency comes from the fact that we implicitly use the Monte Carlo tracking efficiency for taggable tracks. As these tracks can be in dense environments in or near jets, we expect the efficiency to be somewhat less than for isolated tracks. Studies done by embedding Monte Carlo tracks in jets in both data and Monte Carlo events indicate that the Monte Carlo tracking efficiency in dense environments is a few percent higher than in data. We assign a $`\pm `$5% systematic uncertainty to the tagging efficiency for this effect. As described in Section V.2 the systematic uncertainty due to the modeling of the isolation dependence of the tagging efficiency is $`+`$0%, $``$8%. Finally, the statistical uncertainty on the measurement of the SLT tagging efficiency in $`t\overline{t}`$ events, differences between PYTHIA and HERWIG, the uncertainty on the semileptonic branching fraction for B mesons and the estimation of the heavy flavor content of the mistag matrix also contribute to the systematic uncertainties. Adding these contributions in quadrature gives an overall systematic uncertainty for the tagging efficiency of $`+`$8%, $``$11%. Note that the uncertainty on the tagging efficiency affects also the backgrounds determination. The reason is that $`t\overline{t}`$ events need to be subtracted from the pretag sample which is used in Equation 3 to determine the $`W`$+jets background. We take this effect into account when calculating the uncertainty on the cross section.
Uncertainties on the tag matrix are determined by the level of agreement between observed tags and predictions in a variety of samples, as described in Section IV. The uncertainty on the $`W`$+fakes and $`Wb\overline{b}+Wc\overline{c}+Wc`$ prediction is $`\pm `$10%.
To determine the uncertainties on the QCD background prediction in events with three or more jets, we define a control sample from the $`E/_T`$ vs. lepton isolation plane, where the primary lepton isolation parameter $`I`$ is between 0.1 and 0.2 and the event has $`E/_T>20\mathrm{GeV}`$. After subtracting expected contributions from $`W`$ and $`t\overline{t}`$ events, all events in this region are expected to be QCD. We determine the systematic uncertainty on the QCD background using the ratio of the observed over predicted number of events in this control region, which should be 1.0. In the sample where the primary lepton is a muon, we measure 0.5$`\pm `$0.4. In the sample where the primary lepton is an electron, we measure 0.8$`\pm `$0.2. A 50% systematic uncertainty is assigned to the F<sub>QCD</sub> measurement for muons and 20% for electrons. We combine this with the statistical uncertainty on F<sub>QCD</sub>, the uncertainty on the correction factor $`k`$, both given in Table 2 and the 10% systematic uncertainty due to the application of the tag matrix. The total QCD background uncertainty is $`\pm `$67% and $`\pm `$19% for muons and electrons, respectively. These values are determined taking into account the correlation between the estimate of the QCD background and the estimate of the $`W`$+fakes and $`W`$+heavy flavor backgrounds (Equations 4 and 6). We add in quadrature the separate effects on the cross section of the QCD uncertainties for electrons and muons.
The systematic uncertainty on the small Drell-Yan background is dominated by its statistical uncertainty. We also include a 33% relative uncertainty to account for changes in the shape of the Drell-Yan spectrum with the number of jets in the event, as described in Section IV.3. Uncertainties on the Monte Carlo background predictions come from uncertainties in the cross sections for the various processes and from the event sizes of the Monte Carlo samples.
The systematic uncertainties and the corresponding shift of the measured cross section value are summarized in Table 6.
## VII Results
Table 1 shows a summary of the background estimates and the number of SLT tagged events as a function of the number of jets. A “tagged event” is an event with at least one tagged jet. The total background and the $`t\overline{t}`$ expectation are also listed. The line labeled “Corrected Background” corresponds to the background after correcting for the $`t\overline{t}`$ content of the pretag sample, as described in Section IV.1.
We calculate the cross section as:
$$\sigma _{t\overline{t}}=\frac{N_{\mathrm{obs}}N_{\mathrm{bgnd}}}{A_{t\overline{t}}𝑑t},$$
(9)
where $`N_{\mathrm{obs}}`$ is the number of events with $`3`$ jets that are tagged with at least 1 SLT, $`N_{\mathrm{bgnd}}`$ is the corrected background and $`A_{t\overline{t}}`$ is the total acceptance (geometrical acceptance times kinematic acceptance times tagging efficiency), taken from Tables 4 and 5. For events with three or more jets, the total denominator is 1.98 $`\pm `$ 0.28 pb<sup>-1</sup>.
From the number of candidate events with three or more jets, we find a total $`t\overline{t}`$ production cross section of
$`\sigma (p\overline{p}t\overline{t})=5.3\pm 3.3{}_{1.0}{}^{+1.3}\mathrm{pb},`$
where the first uncertainty is statistical and the second is systematic. This cross section value uses acceptances and tagging efficiencies appropriate for a top mass of $`175\mathrm{GeV}/c^2`$. The acceptances and efficiencies, and therefore the calculated cross section, change slightly for other assumed top masses. The calculated cross section is 1% higher assuming a top mass of $`170\mathrm{GeV}/c^2`$, and 5% lower assuming a top mass of $`180\mathrm{GeV}/c^2`$.
Figure 14 shows the number of tags in $`W+1,2,3,4`$ jet events together with the histograms representing the total corrected background with and without the $`t\overline{t}`$ signal expectation, based on the theoretical cross section of $`6.7`$pb for $`M_{\mathrm{top}}=175\mathrm{GeV}/c^2`$.
We examine a number of kinematic distributions of the tagged events and compare with expectations based on the measured signal plus background. Figure 15 shows the $`E_\mathrm{T}`$ distribution of the tagged jets in $`W`$+1 and 2 jets (combined), and in the signal region of $`W`$ plus three or more jets. The $`W`$+1 and 2 jet data-Monte Carlo comparison has a Kolmogorov-Smirnov test (KS) probability of 41%, and the three or more jet comparison has a KS probability of 82%.
Figure 16 compares the $`P_\mathrm{T}`$ distribution of muons identified as SLT tags with expectations from $`t\overline{t}`$ plus backgrounds. The KS probabilities are 6% for $`W`$+1 and 2 jet comparison and 5% for the three or more jet comparison.
Finally, Figure 17 shows the impact parameter significance, defined as the impact parameter divided by its uncertainty, for the SLT tracks and the expectation from signal plus background. The sign of the impact parameter is defined according to whether the track trajectory crosses the jet axis in front of or behind the event primary vertex. The long-lived component from semi-leptonic $`b`$-hadron decays is readily apparent in the shape of the positive impact parameter distribution. The KS probabilities are 12% for $`W`$+1 and 2 jet comparison and 23% for the three or more jet comparison. (Note that Figures 15, 16 and 17 contain 21 entries since one of the events has two jets tagged with SLT.)
## VIII Conclusions
We have measured the total cross section for $`t\overline{t}`$ production through the decay of top pairs into an electron or muon plus multiple jets. We separate signal from background by identifying semileptonic decays of $`b`$ hadrons into muons. The measured $`t\overline{t}`$ production cross section is $`5.3\pm 3.3`$$`{}_{1.0}{}^{}{}_{}{}^{+1.3}`$ pb, consistent with the expectation of 6.7 pb for standard model production and decay of top quark pairs with a mass of $`175\mathrm{GeV}/c^2`$. Distributions of Jet $`E_\mathrm{T}`$ and impact parameter significance for the tagged events, and the distributions of the $`P_\mathrm{T}`$ of the tags, are also consistent with standard model expectations.
The sensitivity of this analysis to test non-standard model $`t\overline{t}`$ production or decay mechanisms is limited by the statistical uncertainty. The combination of this measurement with other CDF II measurements combXsec will yield a significantly more precise value. Future measurements with the full Run II dataset of $`48`$ fb<sup>-1</sup> will provide further factors of approximately four to six in statistical precision. At the same time, significantly larger datasets will provide avenues for reduction of the systematic uncertainties through such things as better understanding of the tag rate in $`W`$+jets events and direct measurement of the tagging efficiency for semileptonic $`b`$-hadron decays.
## IX Acknowledgments
We are grateful to Tim Stelzer for help with studies of the heavy flavor content of $`\gamma `$+jets and $`W`$+jets events using MADEVENT. We thank the Fermilab staff and the technical staffs of the participating institutions for their vital contributions. This work was supported by the U.S. Department of Energy and National Science Foundation; the Italian Istituto Nazionale di Fisica Nucleare; the Ministry of Education, Culture, Sports, Science and Technology of Japan; the Natural Sciences and Engineering Research Council of Canada; the National Science Council of the Republic of China; the Swiss National Science Foundation; the A.P. Sloan Foundation; the Bundesministerium für Bildung und Forschung, Germany; the Korean Science and Engineering Foundation and the Korean Research Foundation; the Particle Physics and Astronomy Research Council and the Royal Society, UK; the Russian Foundation for Basic Research; the Comision Interministerial de Ciencia y Tecnologia, Spain; and in part by the European Community’s Human Potential Programme under contract HPRN-CT-2002-00292, Probe for New Physics. |
warning/0506/cond-mat0506740.html | ar5iv | text | # Collective Excitations in Quantum Hall Liquid Crystals: Single-Mode Approximation Calculations
## I Introduction
For more than two decades two-dimensional electron systems (2DES), and in particular the quantum Hall effect (QHE) qhe ; laughlin83 ; perspectives have been constant sources of complex and unexpected behavior, perhaps with no equal in the realms of condensed matter physics. The unique combination of extremely high mobilities ($`\mu 10^7`$ m/Vs) in GaAs/Al<sub>x</sub>Ga<sub>1-x</sub>As heterostructures, low temperatures ($`T<100`$ mk), enhancement of interactions due to the reduced dimensionality, and relative quenching of the kinetic energy in strong magnetic fields due to Landau level (LL) quantization, has allowed the emergence of complex and striking behavior due to subtle correlation effects. In fact, from the discovery of the integer and fractional QHE’s in the early eighties qhe (leading to the award of two Nobel prizes), to the existence of novel fractionally charged laughlin83 and composite particles cf00 and many newer interesting many-body phenomena, the QHE has been a consistently active and exciting area of research.
The physics of 2DES in partially filled LL’s is inherently complex due to the high degeneracy of the “unperturbed” ground state (i.e. without the Coulomb interaction). The first successful theoretical approach to this system was proposed shortly after the discovery of the FQHE qhe by Laughlin, who proposed his famous trial-wavefunction laughlin83
$$\mathrm{\Psi }_{1/m}=\underset{i<j}{\overset{N}{}}(z_iz_j)^me^{\frac{1}{4}_{k=1}^N|z_k|^2},$$
(1)
to describe states at filling factor $`\nu =1/m`$, with $`m`$ an odd integer. Here $`z_j=x_j+iy_j`$ is the position of the $`j^{\mathrm{th}}`$ electron in the complex plane and we work in units of the magnetic length $`l_0=[\mathrm{}/eB]^{1/2}`$. The “goodness” of Laughlin state originates in that the nodal hyper-surfaces on which the many-body wavefunction vanishes coincide with the ones in which the particles are in contact. Whereas the vanishing of the wavefunction when two (spin aligned) electrons are in contact is required by Fermi statistics, Laughlin state has multiple nodes at those points, thus reducing the Coulomb repulsion.
Later, Jain proposed cf00 the elegant conceptual framework of the composite fermions (CF) which unifies all hierarchies of integer and fractional QHE’s, along with the intermediate states in between QH plateaus: an even number $`2p`$ of vortices is attached to each electron, also lowering the Coulomb energy, and (in a “mean field” approximation) the “magnetic fluxes” associated with the vortices lead to a reduction of the effective magnetic field: $`B^{}=B2p\mathrm{\Phi }_0n_e`$ ($`\mathrm{\Phi }_0=h/e`$ is the flux quantum and $`n_e`$ is the electron density), resulting in an effective filling factor $`\nu ^{}=\nu /(12p\nu )`$. It is easy to see that for the strongest FQHE states $`\nu ^{}`$ is an integer (e.g. for $`\nu =1/3`$ and $`p=1`$ $`\nu ^{}=1`$), leading to the interpretation of the FQHE of electrons as a simple IQHE of the CF’s; whereas the intermediate regions (e.g. $`\nu =1/2,1/4`$) correspond to $`\nu ^{}\mathrm{}`$, i.e. a vanishing effective magnetic field: $`B^{}=0`$.
Since 1999, magnetotransport experiments have uncovered a variety of surprising results at low temperatures ($`T100`$ mK), for example: extreme anisotropies lilly99andothers99 and apparent competition between different ordered phases competition in the intermediate regions between quantum Hall plateaus at high LL’s, the melting transition between a Wigner crystal (WC) at $`\nu 1/7`$,melting17 and the microwave conductivity evidence of structural phase transitions in partially filled LL’s.lewis04 A large body of evidence corresponding to seemingly distinct phenomena may be partly understood in terms of a simple conceptual point: that the many-body states involved have an intrinsic crystalline or liquid crystalline order,foglermoessner ; fradkin99 ; wexlerKT ; brs ; liqcryst04 ; many ; radzihovsky02fogler04 ; oursother ; elasticity be it smectic,foglermoessner ; wexlerKT nematic,fradkin99 ; wexlerKT ; brs ; liqcryst04 tetratic or hexatic.liqcryst04
Generalizations of Laughlin wavefunction \[Eq. (1)\] with discrete broken rotational symmetry (BRS) have been proposed in the past brs ; liqcryst04 ; joynt96balents96 as candidates for nematic or hexatic states brs ; liqcryst04 in order to understand anisotropic transport observed in the intermediate regions,lilly99andothers99 or the melting of the WC at $`\nu =1/7`$.melting17 In fact, the motivation for these states arises from the fact that it is generally expected that melting in 2D may occur through a topological Koszterlitz-Thouless-type (KT) transition.kosterlitz73 For 2D melting, the reliable Kosterlitz-Thouless-Halperin-Nelson-Young (KTHNY) theory kosterlitz73 ; hny predicts that, in fact, an intermediate liquid crystalline phase will exist between a solid and a liquid phase, which will exhibit no translational order, and only a quasi-long-range order for the orientational order below the KT transition temperature. These arguments were used by Wexler and Dorsey wexlerKT to calculate qualitatively correct anisotropic-isotropic transition temperatures for the quantum Hall liquid crystal in the transitional regions at high LL’s.lilly99andothers99
In this paper we consider the spectrum of collective excitations for a family of liquid crystalline states in a partially filled LL. These states are generated so as to satisfy the following criteria which we consider reasonable for understanding the dynamics: (i) states must obey Fermi statistics, i.e., they must be odd under the exchange of any pair of electrons; (ii) the states must be translationally invariant (far enough from boundaries); (iii) there must be a broken rotational symmetry belonging to the proper symmetry group (i.e., $`C_2`$ for a nematic, $`C_4`$ for a tetratic, and $`C_6`$ for a hexatic; additional symmetries are possible in principle, e.g. with a $`C_{10}`$ symmetry, we have not explored such possibilities); (iv) states and excitations must reside entirely in the LLL to avoid the large cyclotron energy cost $`\mathrm{}\omega _c`$. Note, as we will show later, that the properties of any excited LL may be readily obtained from the properties of the LLL.
States that satisfy the abovementioned requirements have been proposed and studied in detail brs ; liqcryst04 ; joynt96balents96 for filling factors $`\nu =1/3`$, 1/5, and 1/7. These are found by splitting the “extra” vortices of the Laughlin (or other CF) states around the electron, while obeying the required symmetries:
$`\mathrm{\Psi }_{1/(2p+1)}^\alpha `$ $`=`$ $`\left\{{\displaystyle \underset{i<j}{\overset{N}{}}}\left[{\displaystyle \underset{\mu =1}{\overset{2p}{}}}(z_iz_j\alpha _\mu )\right]\right\}\times `$ (2)
$`\times {\displaystyle \underset{i<j}{\overset{N}{}}}(z_iz_j)e^{\frac{1}{4}_{k=1}^N|z_k|^2},`$
where the complex numbers $`\alpha _\mu `$ are distributed in pairs of opposite value in the complex plane (to satisfy Fermi statistics). For the states with the highest discrete symmetry at each filling factor we may take
$$\alpha _\mu =\alpha e^{i\mathrm{\hspace{0.17em}2}\pi (\mu 1)/2p},\mu \{1,2,\mathrm{},2p\},$$
(3)
and without loss of generality $`\alpha `$ can be taken to be real. The wavefunction in Eq. (2) represents a homogeneous liquid crystalline state at $`\nu =1/(2p+1)`$, is anti-symmetric, lies entirely in the LLL, and is smoothly connected to the isotropic Laughlin state for $`\alpha =0`$. Figure 1 depicts the nodal distribution for various states of Eq. (2). A remarkable feature of these states is that they posses an underlying charge density wave (CDW), but these CDWs are melted by fluctuations, and overall the system is translationally invariant.liqcryst04
## II The Single Mode Approximation
To calculate the excitation spectrum we use the single mode approximation (SMA),feynman ; gmp-sma which reliably provides the first moment (mean) of the energy of the excitations (for a given wavevector $`𝐤`$) that are coupled to the ground state by means of the density operator.feynman ; gmp-sma ; park00 The SMA was first used by Feynman in 1953 to accurately calculate the spectrum of phonons in superfluid <sup>4</sup>He.feynman The essence of the method originates on the assumption that the ground state of a system of bosons has a scarcity of long-wavelength excitations. Under those circumstances, the variational wavefunction for an excitation corresponding to a density-wave can be written as
$$\varphi _𝐤(𝐫_1,\mathrm{},𝐫_N)=N^{1/2}\rho _𝐤\psi _0(𝐫_1,\mathrm{},𝐫_N),$$
(4)
where $`\rho _𝐤=_{j=1}^Ne^{i𝐤𝐫_j}`$ is the density operator, and $`\psi _0`$ is the many-body ground state (which is, in fact, unknown for <sup>4</sup>He). Note that this trial state automatically builds in the favorable correlations of the ground state. The energy of this excited state, $`\mathrm{\Delta }(𝐤)=\varphi _𝐤|HE_0|\varphi _𝐤/\varphi _𝐤|\varphi _𝐤`$, can be simply evaluated:
$$\mathrm{\Delta }(𝐤)=\frac{N^1\psi _0|\rho _𝐤^{}[H,\rho _𝐤]|\psi _0}{N^1\psi _0|\rho _𝐤^{}\rho _𝐤|\psi _0}\frac{f(𝐤)}{S(𝐤)},$$
(5)
In the last term the numerator is the “oscillator strength” and takes on the universal value $`f(𝐤)=\mathrm{}^2k^2/2m`$, and $`S(𝐤)`$ is the static structure factor, which is directly measurable by means of neutron scattering (it is here, in fact, where the He-He correlations of the ground state are “encoded”). Using the experimental results for $`S(𝐤)`$, Feynman could calculate a spectrum of remarkable quality, showing the phonon-like spectrum at small wavevectors, and a roton minimum at wave-vectors comparable to the inverse interatomic distance.feynman
## III The Single Mode Approximation in the Quantum Hall Effect
The applicability of the SMA to fermion systems is also well established for two- and three-dimensional systems in the absence of magnetic fields, giving a good approximation for the plasmons at long wavelength. For 2DES in presence of a magnetic field it correctly gives the zero-wavevector magnetoplasmon at $`\omega _c=eB/m_e`$, a result that is guaranteed by Kohn’s theorem,kohn61 which states that the dipolar excitation is saturated by the cyclotron mode (this results in the modes of interest—the intra-LL excitations—having quadrupolar character, i.e. with an oscillator strength $`\overline{f}q^4`$).
For excitations fully contained within a single LL, the cyclotron mode is not of primary interest. In 1985 Girvin, MacDonald an Platzman (GMP) proposed an ingenious ansatz for projected excited states:gmp-sma
$$|\psi _𝐪=\overline{\rho }_𝐪|\psi _0,$$
(6)
where $`\overline{\rho }_𝐪`$ is the projected density operator:girvin84
$`\overline{\rho }_𝐪`$ $`=`$ $`{\displaystyle \underset{m,m^{}}{}}0,m^{}|e^{i𝐪𝐫}|0,ma_{0,m^{}}^{}a_{0,m}`$
$`=`$ $`{\displaystyle \underset{j=1}{\overset{N}{}}}\overline{e^{i𝐪𝐫}}={\displaystyle \underset{j=1}{\overset{N}{}}}e^{|q|^2/2}e^{iq^{}z_j/2}e^{iq\frac{}{z_j}},`$
where $`|0,m`$ correspond to single-particle states in the lowest LL and angular momentum $`m`$, and $`a_{0,m}^{}`$ is the particle creator operator for such state. As in Feynman’s ansatz \[Eq. (4)\], Eq. (6) preserves the favorable correlations of the ground state. The exclusion of inter-LL excitations eliminates the problem with the saturation of the dipolar mode. The excited states have a compelling description, in first quantized form:girvin84
$`\overline{\rho }_𝐪\psi (z_1,\mathrm{},z_N)`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{N}{}}}e^{|q|^2/2}e^{iq^{}z_j/2}\times `$
$`\times \psi (z_1,\mathrm{},z_{j1},z_jiq,z_{j+1},\mathrm{},z_N),`$
which corresponds to shifting each electron by $`\widehat{e}_z\times 𝐪`$ and superimposing these $`N`$ configurations with an amplitude $`e^{iq^{}z_j/2}`$.
Similarly to Eq. (5), the excitation spectrum can be readily obtained
$$\overline{\mathrm{\Delta }}_𝐪=\frac{(2N)^1\psi _0|[\overline{\rho }_𝐪^{},[\overline{H},\overline{\rho }_𝐪]]|\psi _0}{N^1\psi _0|\overline{\rho }_𝐪^{}\overline{\rho }_𝐪|\psi _0}\frac{\overline{f}(𝐪)}{\overline{S}(𝐪)}.$$
The projected oscillator strength comes from the non-commutation of the projected density operator with terms in the potential energy part of the Hamiltonian also projected onto the LLL:
$$\overline{H}=\frac{1}{2}\underset{𝐪}{}v_𝐪(\overline{\rho }_𝐪^{}\overline{\rho }_𝐪Ne^{q^2l^2/2}).$$
(9)
Since $`[\overline{\rho }_𝐤,\overline{\rho }_𝐪]=(e^{k^{}q/2}e^{kq^{}/2})\overline{\rho }_{𝐤+𝐪}`$, we find gmp-sma
$`\overline{f}(𝐪)`$ $`=`$ $`2e^{|q|^2/2}{\displaystyle \underset{𝐤}{}}\mathrm{sin}^2({\displaystyle \frac{𝐪\times 𝐤}{2}})\overline{S}(𝐤)\times `$
$`\times [v_{𝐤𝐪}e^{𝐤𝐪}e^{|q|^2/2}v_𝐤],`$
For its part, the projected static structure factor $`\overline{S}(𝐪)`$ can be calculated from:gmp-sma ; girvin84
$$\overline{S}(𝐪)=S(𝐪)(1e^{q^2/2}),$$
(11)
where $`S(𝐪)`$ is the unprojected static structure factor:
$$S(𝐪)1=\rho _0d^2re^{i𝐪𝐫}[g(𝐫)1],$$
(12)
the Fourier transform of the pair correlation function
$$g(𝐫)=n_e^2\underset{ij}{\overset{N}{}}\delta (𝐫_i𝐫)\delta (𝐫_j)$$
(13)
which is obtainable from the ground state via, e.g., Monte Carlo (MC) simulations.gmp-sma ; brs ; liqcryst04 In our case, we considered BRS states \[Eq. (2)\] corresponding to a $`\nu =1/3`$ nematic, a $`\nu =1/5`$ tetratic, and a $`\nu =1/7`$ hexatic. Figure 2 depicts the pair correlation function for various states. In all cases, correlation functions and SMA excitation spectra were computed for numerous $`\alpha `$’s. The angle dependence significantly increases the burden in the MC simulations since the full angle-dependent $`g(𝐫)`$ is needed, rather than the considerably simpler angle-averaged $`g(r)`$ of the isotropic systems. The accurate calculation of $`\overline{f}(𝐪)`$, with its angular-dependent exponentially large factors required high quality $`\overline{S}(𝐪)`$ and hence $`g(𝐫)`$. To put things in perspective, to achieve ca. 1% accuracy in $`g(𝐫)`$, $`𝒪[10,000]`$ counts were accumulated for each $`0.01\times 0.01`$ (in units of $`l_0`$) $`𝐫`$-box in the original histograms (the small boxes were necessary to have precision in the Fourier transform at high wavevectors). This required, for each filling factor $`\nu `$ and anisotropy generating $`\alpha `$, to run approximately 4–8$`\times 10^7`$ MC steps in systems of $`N_e=`$ 200–400 electrons taking 100–200 cpu$`\times `$days of computation for each $`(\nu ,\alpha )`$ in a 2 GHz Athlon cluster (see Ref. \[liqcryst04, \] for further details). A relatively large $`N_e`$ is required so that the simulations are able to reproduce a system in the thermodynamic limit. The fact that $`g(𝐫)1`$ over a large area is a guarantee of that achievement (Fig. 2).
Figure 3 shows the projected static structure factors $`\overline{S}(𝐪)`$ for a $`\nu =1/3`$ nematic, a $`\nu =1/5`$ tetratic,s and a $`\nu =1/7`$ hexatic. From $`\overline{S}(𝐪)`$, the oscillator strength $`\overline{f}(𝐪)`$ is computed using Eq. (III). Analysis of $`\overline{S}(𝐪)`$ and $`\overline{f}(𝐪)`$ shows that both are $`𝒪[q^4]`$ for small $`𝐪`$.gmp-sma This restriction on the the small wavevector behavior originates in Kohn’s theorem,kohn61 as all of the $`𝒪[q^2]`$ pieces in the unprojected parts are saturated by the (uninteresting) inter-LL excitations at $`\mathrm{}\omega _c`$.
Figure 4 presents some of our results for the excitation spectra in the lowest LL. The results are consistent with those obtained by GMP for the isotropic $`\nu =1/3`$ and 1/5 FQHE cases gmp-sma which were qualitatively confirmed experimentally.pinczuk They show that the collective excitation spectrum remains gaped, albeit with a deep magnetoroton, for modes coupled to the ground state via the density operator. Not surprisingly, BRS states have significant anisotropy in their spectra. However, it is noteworthy that for the nematic and tetratic cases the spectrum is singular, with an angle dependence on the excitation energy $`\mathrm{\Delta }(𝐪)`$ as $`𝐪0`$. By contrast, the hexatic liquid crystal has a regular spectrum in the long wavelength limit. The apparent disparity have, of course, to do with the different rotational symmetries of the different states: as $`\overline{\mathrm{\Delta }}=\overline{f}/\overline{S}`$, and both numerator and denominator are $`𝒪[q^4]`$ for small $`𝐪`$, there is no possibility of generating a $`C_6`$ symmetric form with terms that can only depend on $`q_x^4`$, $`q_y^4`$, and $`q_x^2q_y^2`$. This can also be understood from the point of view of an effective elasticity theory for 2DES (valid in the long wavelength limit, $`q0`$):vignale98 the elasticity tensor being of the fourth rank is not compatible with a six-fold rotational symmetry. This elastic interpretation will be published elsewhere.elasticity
The presence of this singular spectrum suggests that microwave conductivity experiments lewis04 ; pinczuk may be able to discern such structures as a signature of, e.g., the quantum Hall nematic suggested by magnetotransport experiments.lilly99andothers99 ; competition
Another interesting feature (Fig. 5) is that the magnetoroton minima and the gap at the origin appears to collapse at some finite anisotropy factor $`\alpha `$. The appearance of elastic modes that go soft may be a precursor of the appearance of charge density waves in the system.liqcryst04 Spectra in higher LL’s can be obtained by multiplying the projected density operator $`\overline{\rho }_𝐪`$ by $`L_L(q^2/2)`$, where $`L_L(x)`$ is the Laguerre polynomial and $`L`$ corresponds to the desired LL.higherLLsVq The rightmost panel of Fig. 4 shows the modification of the SMA spectrum for the first excited LL. Numerical error due to the large $`L_L(q^2/2)`$ factors makes it difficult to get reliable results for higher LL’s at this point.
Theoretical predictions of which type of order is lowest in free energy (required to decide which state is favorable at finite temperatures, as shown in Refs. \[competition, ; melting17, \]) are still incomplete, however, as the entropy is likely to be dominated by gapless modes originating, e.g., from the Goldstone modes associated with the spontaneous breaking of the continuous rotational symmetry of the isotropic states. Generalization of the SMA to gapless rotational modes (see e.g. radzihovsky02fogler04, ) will require, however, 3-body operators which demanding considerably higher computing capabilities.foglerprivate
## IV Summary
Summarizing: in connection with the recent variety of experimental evidence supporting liquid crystalline phases in quantum Hall systems,lilly99andothers99 ; competition ; melting17 ; lewis04 we have calculated the collective excitation spectrum in the SMA approximation for a quantum Hall liquid crystal. We found that the spectrum of excitations coupled to the ground state by the density operator remains gapped, but develops a significant anisotropy, which in the case of the nematic and tetratic liquid crystals has a singular gap in the long wavelength limit.
## V Acknowledgments
We would like to acknowledge helpful discussions with G. Vignale, O. Ciftja, A.T. Dorsey, M. Fogler, A. MacDonald, S. Girvin, E. Fradkin, H. Fertig, and J. Jain. Acknowledgment is made to the University of Missouri Research Board and Research Council, and to the Donors of the Petroleum Research Fund, administered by the American Chemical Society, for support of this research. |
warning/0506/physics0506010.html | ar5iv | text | # NUHEP Report 1010 May 2, 2005 Sifting data in the real world
## 1 Introduction
In an idealized world where all of the data follow a normal (Gaussian) distribution, the use of the $`\chi ^2`$ likelihood technique, through minimization of $`\chi ^2`$, described in detail in A.2, offers a powerful statistical analysis tool when fitting models to a data sample. It allows the phenomenologist to conclude either:
* The model is accepted, based on the value of its $`\chi _{\mathrm{min}}^2`$. It certainly fits well when $`\chi _{\mathrm{min}}^2`$, when compared to $`\nu `$, the numbers of degrees of freedom, has a reasonably high probability ($`\chi _{\mathrm{min}}^2\nu `$). On the other hand, it might be accepted with a much poorer $`\chi _{\mathrm{min}}^2`$, depending on the phenomenologist’s judgment. In any event, the goodness-of-fit of the data to the model is known and an informed judgment can be made.
* Its parameter errors are such that a change of $`\mathrm{\Delta }\chi ^2=1`$ from $`\chi _{\mathrm{min}}^2`$ corresponds to changing a parameter by its standard error $`\sigma `$. These errors and their correlations are summarized in the standard covariance matrix $`C`$ discussed in Appendix A.2.
or
* The model is rejected, because the probability that the data set fits the model is too low, i.e., $`\chi _{\mathrm{min}}^2>>\nu `$.
This decision-making capability (of accepting or rejecting the model) is of primary importance, as is the ability to estimate the parameter errors and their correlations.
Unfortunately, in the real world, experimental data sets are at best only approximately Gaussian and often are riddled with outliers—points far off from a best fit curve to the data, being many standard deviations away. This can be due to many sources, copying errors, bad measurements, wrong calibrations, misassignment of experimental errors, etc. It is this world that our note wishes to address—a world with many data points, and perhaps, many different experiments from many different experimenters, with possibly a significant number of outliers.
In Section 2 we will propose our “Sieve” algorithm, an adaptive technique for discarding outliers while retaining the vast majority of the good data. This then allows us to estimate the goodness-of fit and make a robust determination of both the parameters and their errors—for a discussion of the term “robust”, see Appendix A. In essence, we then retain all of the statistical benefits of the conventional $`\chi ^2`$ technique.
In Sections 3.7.1 and 3.7.2 we will apply the algorithm to high energy $`\overline{p}p`$ and $`pp`$ scattering, as well as to $`\pi ^{}p`$ and $`\pi ^+p`$ scattering. Eight examples of real world experimental data, for both $`\overline{p}p`$ and $`pp`$ scattering and $`\pi ^+p`$ and $`\pi ^{}p`$ scattering, are taken from the Particle Data Group archives and are illustrated in Figures 1, 2, 3 and 4, respectively. The data in Fig. 1 are all of the known published data for the total cross sections $`\sigma _{\overline{p}p}`$ and $`\sigma _{pp}`$ for cms (center of mass) energies greater than 6 GeV. The measured $`\rho _{\overline{p}p}`$ and $`\rho _{pp}`$, where $`\rho `$ is the ratio of the real to the imaginary portion of the forward scattering amplitude, are shown in Fig. 2, again for cms energies greater that 6 GeV. The data in Fig. 3 are all of the known published data for the total cross sections $`\sigma _{\pi ^{}p}`$ and $`\sigma _{\pi ^+p}`$ for cms energies greater that 6 GeV. The measured $`\rho _{\pi ^{}p}`$ and $`\rho _{\pi ^+p}`$ are shown in Fig. 4, again for cms energies greater that 6 GeV. Detailed examination of Figures 1, 2, 3 and 4 show many points far off of the common trend, often at the same energy. Attempts to use the $`\chi ^2`$ technique to fit these data with a model will always come up short. These fits will always return a huge value of $`\chi _{\mathrm{min}}^2/\nu `$, together with model parameters that are likely to be unreliable.
In Section 3, we make three types of computer simulations, generating data normally distributed about a straight line, a constant, and about a parabola, along with outliers—artificial worlds where we know all of the answers, i.e., which points are signal and which are noise. Examples for the straight line, a constant (two cases) and the parabola are shown in Fig. 5, 6, 7 and 8. Details are given in Sections 3.1, 3.3 and 3.6. The noise points in Fig. 5a, 6a, 7a, and Fig. 8a are the diamonds, whereas the signal points are the circles.
The dashed curve in Fig. 5b is the result of a $`\chi ^2`$ fit to all of the noisy data (100 signal plus 20 noise points) in Figure 5a and is not a very good fit to the data. The solid line is the fit with the “ Sieve” algorithm proposed in the next Section. It reproduces nicely the theoretical straight line $`y=12x`$ that was used to computer-generate data that were normally distributed about it, using random numbers. In this case, the 20 noise points penetrated the signal down to a level $`\mathrm{\Delta }\chi _i^2>6`$.
In Fig. 6b we show the results for fitting the constant $`y=10`$. The noise points (diamonds) in Fig. 6a penetrate the signal down to $`\mathrm{\Delta }\chi _i^2>4`$.
In Fig. 7b we show the results for fitting the constant $`y=10`$, where the noise points (diamonds) in Fig. 7a penetrate the signal down to $`\mathrm{\Delta }\chi _i^2>9`$.
In Fig. 8a the data were generated about the parabola $`y=1+2x+0.5x^2`$, with background noise. Figure 8b shows the result of sifting the data according to our Sieve algorithm, described below. The noise points that are retained after invoking our algorithm are the diamonds in Fig. 8b and the circles are the signal points that are retained.
In Sections 3.2 and 3.3, we will calibrate the algorithm with extensive computer-generated numerical simulations and test it for stability and accuracy. The lessons learned from these computer simulations of events are summarized in Section 3.4.
Finally, in Appendix A we give mathematical details about fitting data using the robust $`\mathrm{\Lambda }^2`$ (Lorentzian) maximum likelihood estimator that we employ in our “Sieve” algorithm and in particular, $`\mathrm{\Lambda }_0^2`$, which minimizes the rms (root mean square) widths of the parameter distributions, making them essentially the same as the rms distributions of a $`\chi ^2`$ fit. We also discuss fitting data with the more conventional $`\chi ^2`$ maximum likelihood estimator.
## 2 The Adaptive Sieve Algorithm
### 2.1 Major assumptions
Our major assumptions about the experimental data are:
1. The experimental data can be fitted by a model which successfully describes the data.
2. The signal data are Gaussianly distributed, with Gaussian errors.
3. That we have “outliers” only, so that the background consists only of points “far away” from the true signal.
4. The noise data, i.e. the outliers, do not completely swamp the signal data.
### 2.2 Algorithmic steps
We now outline our adaptive Sieve algorithm, consisting of several steps:
1. Make a robust fit (see Appendix A) of all of the data (presumed outliers and all) by minimizing $`\mathrm{\Lambda }_0^2`$, the Lorentzian squared, defined as
$$\mathrm{\Lambda }_0^2(𝜶;𝒙)\underset{i=1}{\overset{N}{}}\mathrm{ln}\left\{1+0.18\mathrm{\Delta }\chi _i^2(x_i;𝜶)\right\},$$
(1)
described in detail in the Appendix A.4. The $`M`$-dimensional parameter space of the fit is given by $`𝜶=\{\alpha _1,\mathrm{},\alpha _M\}`$, $`𝒙=\{x_1,\mathrm{},x_N\}`$ represents the abscissa of the $`N`$ experimental measurements $`𝒚=\{y_1,\mathrm{},y_N\}`$ that are being fit and $`\mathrm{\Delta }\chi _i^2(x_i;𝜶)\left(\frac{y_iy(x_i;𝜶)}{\sigma _i}\right)^2`$, where $`y(x_i;𝜶)`$ is the theoretical value at $`x_i`$ and $`\sigma _i`$ is the experimental error. As discussed in Appendix A.4, minimizing $`\mathrm{\Lambda }_0^2`$ gives the same total $`\chi _{\mathrm{min}}^2_{i=1}^N\mathrm{\Delta }\chi _i^2(x_i;𝜶)`$ from eq. (1) as that found in a $`\chi ^2`$ fit, as well as rms widths (errors) for the parameters—for Gaussianly distributed data—that are almost the same as those found in a $`\chi ^2`$ fit. The quantitative measure of “far away” from the true signal, i.e., point $`i`$ is an outlier corresponding to Assumption (3), is the magnitude of its $`\mathrm{\Delta }\chi _i^2(x_i;𝜶)=\left(\frac{y_iy(x_i;𝜶)}{\sigma _i}\right)^2`$.
If $`\chi _{\mathrm{min}}^2`$ is satisfactory, make a conventional $`\chi ^2`$ fit to get the errors and you are finished. If $`\chi _{\mathrm{min}}^2`$ is not satisfactory, proceed to step 2.
2. Using the above robust $`\mathrm{\Lambda }_0^2`$ fit as the initial estimator for the theoretical curve, evaluate $`\mathrm{\Delta }\chi _i^2(x_i;𝜶)`$, for each of the $`N`$ experimental points.
3. A largest cut, $`\mathrm{\Delta }\chi _i^2(x_i;𝜶)_{\mathrm{max}}`$, must now be selected. For example, we might start the process with $`\mathrm{\Delta }\chi _i^2(x_i;𝜶)_{\mathrm{max}}=9`$. If any of the points have $`\mathrm{\Delta }\chi _i^2(x_i;𝜶)>\mathrm{\Delta }\chi _i^2(x_i;𝜶)_{\mathrm{max}}`$, reject them—they fell through the “Sieve”. The choice of $`\mathrm{\Delta }\chi _i^2(x_i;𝜶)_{\mathrm{max}}`$ is an attempt to pick the largest “Sieve” size (largest $`\mathrm{\Delta }\chi _i^2(x_i;𝜶)_{\mathrm{max}}`$) that rejects all of the outliers, while minimizing the number of signal points rejected.
4. Next, make a conventional $`\chi ^2`$ fit to the sifted set—these data points are the ones that have been retained in the “Sieve”. This fit is used to estimate $`\chi _{\mathrm{min}}^2`$. Since the data set has been truncated by eliminating the points with $`\mathrm{\Delta }\chi _i^2(x_i;𝜶)>\mathrm{\Delta }\chi _i^2(x_i;𝜶)_{\mathrm{max}}`$, we must slightly renormalize the $`\chi _{\mathrm{min}}^2`$ found to take this into account, by the factor $``$. This effect is discussed later in detail in Section 3.4.
If the renormalized $`\chi _{\mathrm{min}}^2`$, i.e., $`\times \chi _{\mathrm{min}}^2`$ is acceptable—in the conventional sense, using the $`\chi ^2`$ distribution probability function—we consider the fit of the data to the model to be satisfactory and proceed to the next step. If the renormalized $`\chi _{\mathrm{min}}^2`$ is not acceptable and $`\mathrm{\Delta }\chi _i^2(x_i;𝜶)_{\mathrm{max}}`$ is not too small, we pick a smaller $`\mathrm{\Delta }\chi _i^2(x_i;𝜶)_{\mathrm{max}}`$ and go back to step 3. The smallest value of $`\mathrm{\Delta }\chi _i^2(x_i;𝜶)_{\mathrm{max}}`$ that makes much sense, in our opinion, is $`\mathrm{\Delta }\chi _i^2(x_i;𝜶)_{\mathrm{max}}>2`$. After all, one of our primary assumptions is that the noise doesn’t swamp the signal. If it does, then we must discard the model—we can do nothing further with this model and data set!
5. From the $`\chi ^2`$ fit that was made to the “sifted” data in the preceding step, evaluate the parameters $`𝜶`$. Next, evaluate the $`M\times M`$ covariance (squared error) matrix of the parameter space which was found in the $`\chi ^2`$ fit. We find the new squared error matrix for the $`\mathrm{\Lambda }^2`$ fit by multiplying the covariance matrix by the square of the factor $`r_{\chi ^2}`$ (for example, as shown later in Section 3.2.2, $`r_{\chi ^2}1.02,1.05`$, 1.11 and 1.14 for $`\mathrm{\Delta }\chi _i^2(x_i;𝜶)_{\mathrm{max}}=9`$, 6, 4 and 2, respectively ). The values of $`r_{\chi ^2}>1`$ reflect the fact that a $`\chi ^2`$ fit to the truncated Gaussian distribution that we obtain—after first making a robust fit—has a rms (root mean square) width which is somewhat greater than the rms width of the $`\chi ^2`$ fit to the same untruncated distribution. Extensive computer simulations, summarized in Section 3.4, demonstrate that this robust method of error estimation yields accurate error estimates and error correlations, even in the presence of large backgrounds.
You are now finished. The initial robust $`\mathrm{\Lambda }_0^2`$ fit has been used to allow the phenomenologist to find a sifted data set. The subsequent application of a $`\chi ^2`$ fit to the sifted set gives stable estimates of the model parameters $`𝜶`$, as well as a goodness-of-fit of the data to the model when $`\chi _{\mathrm{min}}^2`$ is renormalized for the effect of truncation due to the cut $`\mathrm{\Delta }\chi _i^2(x_i;𝜶)_{\mathrm{max}}.`$ Model parameter errors are found when the covariance (squared error) matrix of the $`\chi ^2`$ fit is multiplied by the appropriate factor $`(r_{\chi ^2})^2`$ for the cut $`\mathrm{\Delta }\chi _i^2(x_i;𝜶)_{\mathrm{max}}`$.
It is the combination of using both $`\mathrm{\Lambda }_0^2`$ (robust) fitting and $`\chi ^2`$ fitting techniques on the sifted set that gives the Sieve algorithm its power to make both a robust estimate of the parameters $`𝜶`$ as well as a robust estimate of their errors, along with an estimate of the goodness-of-fit.
Using this same sifted data set, you might then try to fit to a different theoretical model and find $`\chi _{\mathrm{min}}^2`$ for this second model. Now one can compare the probability of each model in a meaningful way, by using the $`\chi ^2`$ probability distribution function of the numbers of degrees of freedom for each of the models. If the second model had a very unlikely $`\chi _{\mathrm{min}}^2`$, it could now be eliminated. In any event, the model maker would now have an objective comparison of the probabilities of the two models.
### 2.3 Evaluating the Sieve algorithm
We will give two separate types of examples which illustrate the Sieve algorithm. In the first type, we computer-generated data, normally distributed about
* a straight line, along with random noise to provide outliers,
* a constant, along with random noise to provide outliers,
* a parabola, with background noise normally distributed about a slightly different parabola,
the details of which are described below. The advantage here, of course, is that we know which points are signal and which points are noise.
For our real world example, we took four types of experimental data for elementary particle scattering from the archives of the Particle Data Group. For all energies above 6 GeV, we took total cross sections and $`\rho `$-values and made a fit to these data. These were all published data points and the entire sample was used in our fit. We then made separate fits to
* $`\overline{p}p`$ and $`pp`$ total cross sections and $`\rho `$-values,
* $`\pi ^{}p`$ and $`\pi ^+p`$ total cross sections $`\sigma `$ and $`\rho `$-values,
using eqns. (7), (8) and (9) below.
## 3 Studies using large computer-generated data sets
Extensive computer simulations were made using the straight line model $`y_i=12x_i`$ and the constant model $`y_i=10`$. Over 500,000 events were computer-generated, with normal distributions of 100 signal points per event, some with no noise and others with 20% and 40% noise added, in order to investigate the accuracy and stability of the “Sieve” algorithm. The cuts $`\mathrm{\Delta }\chi _i^2>9`$, 6, 4 and 2 were investigated in detail.
### 3.1 A straight line model
An event consisted of generating 100 signal points plus either 20 or 40 background points, for a total of 120 or 140 points, depending on the background level desired. Let RND be a random number, uniformly distributed from 0 to 1. Using random number generators, the first 100 points used $`x_i=10\times \mathrm{RND}`$, where $`i`$ is the point number. This gives a signal randomly distributed between $`x=0`$ and $`x=10`$. For each point $`x_i`$, a theoretical value $`\overline{y}_i`$ was found using $`\overline{y}_i=12x_i`$. Next, the value of $`\sigma _i`$, the “experimental error”, i.e, the error bar assigned to point $`i`$, was generated as $`\sigma _i=a_i+\alpha _i\times \mathrm{RND}`$. Using these $`\sigma _i`$, the $`y_i`$’s were generated, normally distributed about the value of $`\overline{y}_i`$ For $`i=1`$ to 50, $`a_i=0.2,\alpha _i=1.5`$, and for $`i=51`$ to 100, $`a_i=0.2,\alpha _i=3`$. This sample of 100 points made up the signal.
The 40 noise points, $`i=101`$ to 140 were generated as follows. Each point was assigned an “experimental error” $`\sigma _i=a_i+\alpha _i\times \mathrm{RND}`$. The $`x_i`$ were generated as $`x_i=d_i+\delta _i\times \mathrm{RND}`$. In order to provide outliers, the value of $`y_i`$ was fixed at $`y_i=12x_i+f_{\mathrm{cut}}\times \mathrm{Sign}_\mathrm{i}\times (b_i+\beta _i)\times \sigma _i`$ and the points were then placed at this fixed value of $`y_i`$ and given the “experimental error” $`\sigma _i`$. The parameter $`f_{\mathrm{cut}}`$ depended only on the value of $`\mathrm{\Delta }\chi _i^2(x_i;𝜶)_{\mathrm{max}}`$ that was chosen, being 1.9, 2.8, 3.4 or 4, for $`\mathrm{\Delta }\chi _i^2(x_i;𝜶)_{\mathrm{max}}=2`$, 4, 6 or 9, respectively, and was independent of $`i`$. These choices of $`f_{\mathrm{cut}}`$ made outliers that only existed for values of $`\mathrm{\Delta }\chi _i^2(x_i;𝜶)>\mathrm{\Delta }\chi _i^2(x_i;𝜶)_{\mathrm{max}}`$.
For $`i=101`$ to 116, $`d_i=0,\delta _i=10,a_i=0.75,\alpha _i=0.5,b_i=1.0,\beta _i=0.6`$. To make “doubles” at the same $`x_i`$ as a signal point, if $`y_{i100}>12x_{i100}`$ we pick Sign$`{}_{i}{}^{}=+1`$; otherwise Sign$`{}_{i}{}^{}=1`$, so that the outlier is on the same side of the reference line $`12x_i`$ as is the signal point.
For $`i=117`$ to 128, $`d_i=0,\delta _i=10,a_i=0.5,\alpha _i=0.5,b_i=1.0,\beta _i=0.6`$; Sign<sub>i</sub> was randomly chosen as +1 or -1. This generates outliers randomly distributed above and below the reference line, with $`x_i`$ randomly distributed from 0 to 10.
For $`i=129`$ to 140, $`d_i=8,\delta _i=2,a_i=0.5,\alpha _i=0.5,b_i=1.0,\beta _i=0.6`$; Sign<sub>i</sub> = +1. This makes points in a “corner” of the plot, since $`x_i`$ is now randomly distributed at the “edge” of the plot, between 8 and 10. Further, all of this points are above the line, since Sign<sub>i</sub> is fixed at +1, giving these points a large lever arm in the fit.
For the events generated with 20 noise points, the above recipes for background were simply halved. An example of such an event containing 120 points, for which $`\mathrm{\Delta }\chi _i^2(x_i;𝜶)_{\mathrm{max}}=6`$, is shown in Fig. 5a, with the 100 squares being the normally distributed data and the 20 circles being the noise data.
After a robust fit to the entire 120 points, the sifted data set retained 100 points after the $`\mathrm{\Delta }\chi _i^2>6`$ condition was applied. This fit had $`\chi _{\mathrm{min}}^2=88.69`$, with an expected $`\chi ^2=\nu =98`$, giving $`\chi _{\mathrm{min}}^2/\nu =0.905`$. Using a renormalization factor $`=1/0.901`$, we get a renormalized $`\chi _{\mathrm{min}}^2/\nu =1.01`$—see Section 3.4 for details of the renormalization factor. After using the Sieve algorithm, by minimizing $`\chi ^2`$ for the sifted set, we found that the best-fit straight line, $`y=<a>+<b>x`$, had $`<a>=0.998\pm 0.12`$ and $`<b>=2.014\pm 0.020`$. The parameter errors given above come from multiplying the errors found in a conventional $`\chi ^2`$ fit to the sifted data by the factor $`r_{\chi ^2}=1.05`$—for details see Section 3.4. This turns out to be a high probability fit with a probability of 0.48 (since the renormalized $`\chi _{\mathrm{min}}^2/\nu =1.01`$, whereas we expect $`<\chi ^2/\nu >=1.0\pm 0.14`$).
Figure 5b shows the results after the use of the Sieve procedure with $`\mathrm{\Delta }\chi _{i}^{2}{}_{\mathrm{max}}{}^{}=6`$. Of the original 120 points, all 100 of the signal points were retained (squares), while no noise points (diamonds) were retained. The solid line is the best $`\chi ^2`$ fit, $`y=0.9982.014x`$.
Had we applied a $`\chi ^2`$ minimization to original 120 point data set, we would have found $`\chi ^2=570`$, which has infinitesimal statistical probability. The straight line resulting from that fit, $`y=0.9251.98x`$, is also shown in Fig. 5b as the dot-dashed curve. For large x, it tends to overestimate the true values.
To investigate the stability of our procedure with respect to our choice of $`\mathrm{\Delta }\chi _i^2`$, we reanalyzed the full data set for the cut-off, $`\mathrm{\Delta }\chi _{i}^{2}{}_{\mathrm{max}}{}^{}=4`$. The evaluation of the parameters $`a`$ and $`b`$ was completely stable, essentially independent of the choice of $`\mathrm{\Delta }\chi _i^2`$. The robustness of this procedure on this particular data set is evident.
### 3.2 Distributional widths for the straight line model
We now generate extensive computer simulations of data sets resulting from the straight line $`y_i=12x_i`$ using the recipe of Section 3.1, with and without outliers, in order to test the Sieve algorithm. We have generated 50,000 events with 20% background and 50,000 events with 40% background, for each cut $`\mathrm{\Delta }\chi _{i}^{2}{}_{\mathrm{max}}{}^{}=9`$, 6, 4 and 2. We also generated 100,000 Gaussianly distributed events with no noise.
#### 3.2.1 Case 1
We generated 100,000 Gaussianly distributed events with no noise. Let $`a`$ and $`b`$ be the intercept and slope of the straight line $`y=12x`$ and define $`<a>`$ as the average $`a`$, $`<b>`$ as the average $`b`$ found for the 100,000 straight-line events, each generated with 100 data points, using both a $`\mathrm{\Lambda }_0^2`$ (robust) fit and a $`\chi ^2`$ fit. The purpose of this exercise was to find $`r(\mathrm{\Lambda }_0^2)`$, the ratio of the $`\mathrm{\Lambda }_0^2`$ rms parameter width $`\sigma (\mathrm{\Lambda }_0^2)`$ divided by $`\mathrm{\Sigma }`$, the parameter error from the $`\chi ^2`$ fit, i.e.
$$r_a(\mathrm{\Lambda }_0^2)\frac{\sigma _a(\mathrm{\Lambda }_0^2)}{\mathrm{\Sigma }_a},r_b(\mathrm{\Lambda }_0^2)\frac{\sigma _b(\mathrm{\Lambda }_0^2)}{\mathrm{\Sigma }_b},$$
as well as demonstrate that there were no biases (offsets) in parameter determinations found in $`\mathrm{\Lambda }^2`$ and $`\chi ^2`$ fits.
The measured offsets $`1<a_{\chi ^2}>`$, $`1<a_{\mathrm{\Lambda }^2}>`$, $`2<b_{\chi ^2}>`$ and $`2<b_{\mathrm{\Lambda }^2}>`$ were all numerically compatible with zero, as expected, indicating that the parameter expectations were not biased.
Let $`\sigma `$ be the rms width of a parameter distribution and $`\mathrm{\Sigma }`$ the error from the $`\chi ^2`$ covariant matrix. We found:
$`\sigma _a(\chi ^2)`$ $`=`$ $`0.139\pm 0.002\mathrm{and}\mathrm{\Sigma }_a=0.138`$
$`\sigma _b(\chi ^2)`$ $`=`$ $`0.0261\pm 0.003\mathrm{and}\mathrm{\Sigma }_b=0.0241,`$
showing that the rms widths $`\sigma `$ and parameter errors $`\mathrm{\Sigma }`$ were the same for the $`\chi ^2`$ fit, as expected. Further, the width ratios $`r`$ for the $`\mathrm{\Lambda }_0^2`$ fit are given by
$`r_a(\mathrm{\Lambda }_0^2)`$ $`=`$ $`1.034\pm 0.010`$
$`r_b(\mathrm{\Lambda }_0^2)`$ $`=`$ $`1.029\pm 0.011,`$
demonstrating that:
* the $`r`$’s of the $`\mathrm{\Lambda }_0^2`$ are almost as good as that of the $`\chi ^2`$ distribution, $`r(\chi ^2)=1`$.
* the ratios of the rms $`\mathrm{\Lambda }^2`$ width to the rms $`\chi ^2`$ width for both parameters $`a`$ and $`b`$ are the same, i.e., we can now simply write
$$r_{\mathrm{\Lambda }^2}=\frac{\sigma _{\mathrm{\Lambda }^2}}{\mathrm{\Sigma }}1.03.$$
(2)
Finally, we find that $`1<\chi ^2/\nu >=0.00034\pm 0.00044`$, which is approximately zero, as expected.
#### 3.2.2 Case 2
For Case 2, we investigate data generated with 20% and 40% noise that have been subjected to the adaptive Sieve algorithm, i.e. the sifted data after cuts of $`\mathrm{\Delta }\chi _i^2(x_i;𝜶)_{\mathrm{max}}=9`$, 6, 4 and 2. We investigated this truncated sample to measure possible biases and to obtain numerical values for $`r`$’s.
We generated 50,000 events, each with 100 points normally distributed and with either 20 or 40 outliers, for each cut. A robust fit was made to the entire sample (either 120 or 140 points) and we sifted the data, rejecting all points with either $`\mathrm{\Delta }\chi _i^2(x_i;𝜶)>9,`$ 6, 4 and 2, according to how the data were generated. A conventional $`\chi ^2`$ analysis was then made to the sifted data. The results are summarized in Table 1. As before, we found that the widths from the $`\chi ^2`$ fit were slightly smaller than the widths from a robust fit, so we adopted only the results for the $`\chi ^2`$ fit.
There were negligible offsets $`1<a>`$ and $`2<b>`$, being $`1`$ to $`5\%`$ of the relevant rms widths, $`\sigma _a`$ and $`\sigma _b`$, for both the robust and $`\chi ^2`$ fits.
In any individual $`\chi ^2`$ fit to the $`j`$th data set, one measures $`a_j,b_i,\mathrm{\Sigma }_{a_j},\mathrm{\Sigma }_{b_j}`$ and $`(\chi _{\mathrm{min}}^2/\nu )_j`$. Thus, we characterize all of our computer simulations in terms of these 7 observables.
We again find that the $`r_{\chi ^2}`$ values—defined as $`\sigma /\mathrm{\Sigma }`$—are the same, whether we are measuring $`a`$ or $`b`$. They are given by $`r_{\chi ^2}=\sigma /\mathrm{\Sigma }=1.034`$, 1.054, 1.098 and 1.162 for the cuts $`\mathrm{\Delta }\chi _i^2(x_i;𝜶)_{\mathrm{max}}=9,`$ 6, 4 and 2, respectively. Further, they are the same for 20% noise and 40% noise, since the cuts rejected all of the noise points. In addition, the $`r`$ values were found to be the same as the $`r`$ values for the case of truncated pure signal, using the same $`\mathrm{\Delta }\chi _i^2(x_i;𝜶)_{\mathrm{max}}`$ cuts. The signal retained was 99.7, 98.57, 95.5 and 84.3 % for the cuts $`\mathrm{\Delta }\chi _i^2(x_i;𝜶)_{\mathrm{max}}=9,`$ 6, 4 and 2, respectively—see Section 3.4 and eq. (6) for theoretical values of the amount of signal retained.
We experimentally determine the rms widths $`\sigma `$ (the errors of the parameter) by multiplying the $`r`$ value, a known quantity independent of the particular event, by the appropriate $`\mathrm{\Sigma }`$ which is measured for that event, i.e.,
$`\sigma _a`$ $`=`$ $`\mathrm{\Sigma }_a\times r_{\chi ^2}`$
$`\sigma _b`$ $`=`$ $`\mathrm{\Sigma }_b\times r_{\chi ^2}.`$
The rms widths are now determined for any particular data set by multiplying the known factors $`r_{\chi ^2}`$ by the appropriate $`\mathrm{\Sigma }`$ found (measured) from the covariant matrix of the $`\chi ^2`$ fit of that data set.
Also shown in Table 1 are the values of $`\chi _{\mathrm{min}}^2/\nu `$ found for the various cuts. We will compare these results later with those for the constant case, in Section 3.3
We again see that a sensible approach for data analysis–even where there are large backgrounds of $`40\%`$—is to use the parameter estimates for $`a`$ and $`b`$ from the truncated $`\chi ^2`$ fit and assign their errors as
$`\sigma _a`$ $`=r_{\chi ^2}\mathrm{\Sigma }_a`$
$`\sigma _b`$ $`=r_{\chi ^2}\mathrm{\Sigma }_b,`$ (3)
where $`r_{\chi ^2}`$ is a function of the $`\mathrm{\Delta }\chi _{i}^{2}{}_{\mathrm{max}}{}^{}`$ cut utilized. Before estimating the goodness-of-fit, we must renormalize the observed $`\chi _{\mathrm{min}}^2/\nu `$ by the appropriate numerical factor for the $`\mathrm{\Delta }\chi _{i}^{2}{}_{\mathrm{max}}{}^{}`$ cut used.
This strategy of using an adaptive $`\mathrm{\Delta }\chi _i^2(x_i;𝜶)_{\mathrm{max}}`$ cut minimizes the error assignments, guarantees robust fit parameters with no significant bias and also returns a goodness-of-fit estimate.
### 3.3 The constant model, $`y_i=10`$
For this case, we investigate a different theoretical model ($`y_i=10`$) with a different background distribution, to measure the values of $`r_{\chi ^2}`$ and $`<\chi _{\mathrm{min}}^2/\nu >`$.
An event consisted of generating 100 signal points plus either 20 or 40 background points, for a total of 120 or 140 points, depending on the background level desired. Again, let RND be a random number, uniformly distributed from 0 to 1. Using random number generators, for the first 100 points $`i`$, a theoretical value $`\overline{y}_i=10`$ was chosen. Next, the value of $`\sigma _i`$, the “experimental error”, i.e, the error bar assigned to point $`i`$, was generated as $`\sigma _i=a_i+\alpha _i\times \mathrm{RND}`$. Using these $`\sigma _i`$, the $`y_i`$’s were generated, normally distributed about the value of $`\overline{y}_i=10`$ . For $`i=1`$ to 50, $`a_i=0.2,\alpha _i=1.5`$, and for $`i=51`$ to 100, $`a_i=0.2,\alpha _i=3`$. This sample of 100 points made up the signal.
The 40 noise points, $`i=101`$ to 140 were generated as follows. Each point was assigned an “experimental error” $`\sigma =a_i+\alpha _i\times \mathrm{RND}`$. In order to provide outliers, the value of $`y_i`$ was fixed at $`y_i=10+f_{\mathrm{cut}}\times \mathrm{sign}_\mathrm{i}\times (b_i+\beta _i)\times \sigma _i`$ and the points were then placed at this fixed value of $`y_i`$ and given the “experimental error” $`\sigma _i`$. The parameter $`f_{\mathrm{cut}}`$ depended only on the value of $`\mathrm{\Delta }\chi _i^2(x_i;𝜶)_{\mathrm{max}}`$ that was chosen, being 1.9, 2.8, 3.4 or 4, for $`\mathrm{\Delta }\chi _i^2(x_i;𝜶)_{\mathrm{max}}=2`$, 4, 6 or 9, respectively, and was independent of $`i`$.
For $`i=101`$ to 116, $`a_i=0.75,\alpha _i=0.5,b_i=1.0,\beta _i=0.6`$; Sign<sub>i</sub> was randomly chosen at +1 or -1.
For $`i=117`$ to 128, $`a_i=0.5,\alpha _i=0.5,b_i=1.0,\beta _i=0.6`$; This generates outliers randomly distributed above and below the reference line, with $`x_i`$ randomly distributed from 0 to 10.
For $`i=129`$ to 140, $`a_i=0.5,\alpha _i=0.5,b_i=1.0,\beta _i=0.6`$; Sign<sub>i</sub> = +1. This forces 12 points to be greater than 10, since Sign<sub>i</sub> is fixed at +1. For the events generated with 20 noise points, the above recipes for background were simply halved.
Two examples of events with 40 background points are shown in Figures 6a and 7a, with the 100 squares being the normally distributed data and the 40 circles being the noise data.
In Fig. 6b we show the results after using the cut $`\mathrm{\Delta }\chi _i^2(x_i;𝜶)_{\mathrm{max}}=4`$. No noise points (diamonds) were retained, and 98 signal points (circles) are shown. The best fit, $`y=9.98\pm 0.074`$, is the solid line, whereas the dashed-dot curve is the fit to all 140 points. The observed $`\chi _{\mathrm{min}}^2/\nu =0.84`$ yields a renormalized value $`\times \chi _{\mathrm{min}}^2/\nu =1.09`$, in good agreement with the expected value $`\chi _{\mathrm{min}}^2/\nu =1\pm 0.14`$. If we had fit to the entire 140 points, we would find $`\chi _{\mathrm{min}}^2/\nu =4.39`$, with the fit being the dashed-dot curve.
In Fig. 7b we show the results after using the cut $`\mathrm{\Delta }\chi _i^2(x_i;𝜶)_{\mathrm{max}}=9`$. No noise points (diamonds) were retained, and 98 signal points (circles) are shown. The best fit, $`y=10.05\pm 0.074`$, is the solid line, whereas the dashed-dot curve is the fit to all 140 points. The observed $`\chi _{\mathrm{min}}^2/\nu =1.08`$ yields a renormalized value $`\times \chi _{\mathrm{min}}^2/\nu =1.11`$, in good agreement with the expected value $`\chi _{\mathrm{min}}^2/\nu =1\pm 0.14`$. If we had fit to the entire 140 points, we would find $`\chi _{\mathrm{min}}^2/\nu =8.10`$, with the fit being the dashed-dot curve. The details of the renormalization of $`\chi _{\mathrm{min}}^2/\nu `$ and the assignment of the errors are given in Section 3.4
We computer-generated a total of 500,000 events, 50,000 events with 20% noise and an additional 50,000 events with 40% noise, for each of the cuts $`\mathrm{\Delta }\chi _i^2>9`$, 6, 4 and 2, and 100,000 events with no noise.
For the sample with no cut and no noise, we found $`r_{\mathrm{\Lambda }_0^2}=1.03\pm 0.02`$, equal to the value $`r_{\mathrm{\Lambda }_0^2}=1.03`$ that was found for the straight line case.
Again, we found that our results for $`r_{\chi ^2}`$ were independent of background, as well as model, and only depended on the cut. We also found that the biases (offsets) for the constant case, $`(10<a_{\chi 2}>)`$, although non-zero for the noise cases, were small in comparison to $`\sigma `$, the rms width.
The results for cuts $`\mathrm{\Delta }\chi _{i}^{2}{}_{\mathrm{max}}{}^{}=9`$, 6, 4 and 2 are detailed in Table 1. We see in Table 1, compared with the straight line results of Section 3.2.2, that the $`r_{\chi ^2}`$ values for the constant case are essentially identical, as expected. Further, we find the same results for the values of $`\chi _{\mathrm{min}}^2/\nu `$ as a function of the cut $`\mathrm{\Delta }\chi _i^2(x_i;𝜶)_{\mathrm{max}}`$.
### 3.4 Lessons learned from computer studies of a straight line model and a constant model
* As found in Sections 3.2.2 and 3.3 and detailed in Table 1, we have universal values of $`r_{\chi ^2}`$ and $`<\chi _{\mathrm{min}}^2>/\nu `$, as a function of the cut $`\mathrm{\Delta }\chi _i^2(x_i;𝜶)_{\mathrm{max}}`$, independent of both background and model.
* A sensible conservative approach for large backgrounds (less than or the order 40%) is to use the parameter estimates from the $`\chi ^2`$ fit to the sifted data and assign the parameter errors to the fitted robust parameters as
$`\sigma (\chi ^2)`$ $`=`$ $`r_{\chi ^2}\times \mathrm{\Sigma },`$
where $`r_{\chi ^2}`$ is a function of the cut $`\mathrm{\Delta }\chi _i^2(x_i;𝜶)_{\mathrm{max}}`$, given by the average of the straight line and constant cases of Table 1. This strategy gives us a minimum parameter error, with only very small biases to the parameter estimates.
* We must then renormalize the value found for $`\chi _{\mathrm{min}}^2/\nu `$ by the appropriate averaged value of $`<\chi _{\mathrm{min}}^2>/\nu `$ for the straight line and constant case, again as a function of the cut $`\mathrm{\Delta }\chi _i^2(x_i;𝜶)_{\mathrm{max}}`$.
* Let us define $`\mathrm{\Delta }`$ as the $`\mathrm{\Delta }\chi _{i}^{2}{}_{\mathrm{max}}{}^{}`$ cut and $``$ as the renormalization factor that multiplies $`\chi _{\mathrm{min}}^2/\nu `$.
We find from inspection of Cases 1 to 2 for the straight line and of Section 3.3 for the case of the constant fit that a best fit parameterization of $`r_{\chi ^2}`$, valid for $`\mathrm{\Delta }2`$ is given by
$$r_{\chi ^2}=1+0.246e^{0.263\mathrm{\Delta }}.$$
(4)
We note that $`^1`$, for large $`\nu `$, is given analytically by
$`^1`$ $``$ $`{\displaystyle _\sqrt{\mathrm{\Delta }}^{+\sqrt{\mathrm{\Delta }}}}x^2e^{x^2/2}𝑑x/{\displaystyle _\sqrt{\mathrm{\Delta }}^{+\sqrt{\mathrm{\Delta }}}}e^{x^2/2}𝑑x`$ (5)
$`=`$ $`1{\displaystyle \frac{2}{\sqrt{\pi }}}{\displaystyle \frac{e^{\mathrm{\Delta }/2}}{\mathrm{erf}(\sqrt{\mathrm{\Delta }/2})}}.`$
Graphical representations of $`r_{\chi ^2}`$ and $`^1`$ are shown in Figures 9a and 9b, respectively. Some numerical values are given in Table 1 and are compared to the computer-generated values found numerically for the straight line and constant cases. The agreement is excellent.
* Let us define $`\sigma _0`$ as the rms parameter width that we would have had for a $`\chi ^2`$ fit to the uncut sample, where the sample had had no background, and define $`\mathrm{\Sigma }_0`$ the error found from the covariant matrix. They are, of course, equal to each other, as well as being the smallest error possible. We note that the ratio $`\sigma /\sigma _0=r_{\chi ^2}\times \mathrm{\Sigma }/\mathrm{\Sigma }_0`$. This ratio is a function of the cut $`\mathrm{\Delta }`$ through both $`r_{\chi ^2}`$ and $`\mathrm{\Sigma }`$, since for a truncated distribution, $`\mathrm{\Sigma }/\mathrm{\Sigma }_0`$ depends inversely on the square root of the fraction of signal points that survive the cut $`\mathrm{\Delta }`$. In particular, the survival fraction $`S.F.`$ is given by
$$S.F.=_\sqrt{\mathrm{\Delta }}^{+\sqrt{\mathrm{\Delta }}}\frac{1}{\sqrt{2\pi }}e^{x^2/2}dx=\mathrm{erf}(\sqrt{\mathrm{\Delta }/2})$$
(6)
and is 99.73, 98.57, 95.45 and 84.27 % for the cuts $`\mathrm{\Delta }=9,`$ 6, 4 and 2, respectively. The survival fraction $`S.F.`$ is shown in Table 1 as a function of the cut $`\mathrm{\Delta }\chi _{i}^{2}{}_{\mathrm{max}}{}^{}`$, as well as is the ratio $`\sigma /\sigma _0`$. We note that the true cost of truncating a Gaussian distribution, i.e., the enlargement of the error due to truncation, is not $`r_{\chi ^2}`$, but rather $`r_{\chi ^2}/\sqrt{S.F.}`$, which ranges from $`1.02`$ to 1.25 when the cut $`\mathrm{\Delta }\chi _{i}^{2}{}_{\mathrm{max}}{}^{}`$ goes from 9 to 2. This rapid loss of accuracy is why the errors become intolerable for cuts $`\mathrm{\Delta }\chi _{i}^{2}{}_{\mathrm{max}}{}^{}`$ smaller than 2.
### 3.5 Fitting strategy
We find that an effective strategy for eliminating noise and making robust parameter estimates, together with robust error assignments, is:
1. Make an initial $`\mathrm{\Lambda }_0^2`$ fit to the entire data sample. If $`\chi _{\mathrm{min}}^2/\nu `$ is satisfactory, then make a standard $`\chi ^2`$ fit to the data and you are finished. If not, then proceed to the next step.
2. Pick a large value of $`\mathrm{\Delta }\chi _i^2(x_i;𝜶)_{\mathrm{max}}`$, e.g., $`\mathrm{\Delta }\chi _i^2(x_i;𝜶)_{\mathrm{max}}=9.`$
3. Obtain a sifted sample by throwing away all points with $`\mathrm{\Delta }\chi _i^2(x_i;𝜶)>\mathrm{\Delta }\chi _i^2(x_i;𝜶)_{\mathrm{max}}`$.
4. Make a conventional $`\chi ^2`$ fit to the sifted sample. For your choice of $`\mathrm{\Delta }\chi _i^2(x_i;𝜶)_{\mathrm{max}}`$, find $`^1`$ from eq. (5). If the renormalized value $`\times \chi _{\mathrm{min}}^2/\nu `$ is sufficiently near 1, i.e., the goodness-of-fit is satisfactory, then go to the next step. If, on the other hand, $`\times \chi _{\mathrm{min}}^2/\nu `$ is too large, pick a smaller value of $`\mathrm{\Delta }\chi _i^2(x_i;𝜶)_{\mathrm{max}}`$ and go to step 3 (for example, if you had used a cut of 9, now pick $`\mathrm{\Delta }\chi _i^2(x_i;𝜶)_{\mathrm{max}}=6`$ and start again). Finally, if you reach $`\mathrm{\Delta }\chi _i^2(x_i;𝜶)_{\mathrm{max}}=2`$ and you still don’t have success, quit—the background has penetrated too much into the signal for the “Sieve” algorithm to work properly.
5. a) Use the parameter estimates found from the $`\mathrm{\Delta }\chi _i^2(x_i;𝜶)_{\mathrm{max}}`$ fit in the previous step.
b) Find a new squared error matrix by multiplying the covariant matrix $`C`$ found in the $`\chi ^2`$ fit by $`(r_{\chi ^2})^2`$. Use the value of $`r_{\chi ^2}`$ found in eq. (4) for the chosen value of the cut $`\mathrm{\Delta }\chi _i^2(x_i;𝜶)_{\mathrm{max}}`$ to obtain a robust error estimate essentially independent of background distribution.
You are now finished. You have made a robust determination of the parameters, their errors and the goodness-of-fit.
The renormalization factors $``$ are only used in estimating the value of the goodness-of-fit, where small changes in this value are not very important. Indeed, it hardly matters if the estimated renormalized $`\chi ^2/\nu `$ is between 1.00 and 1.01—the possible variation of the expected renormalized $`\chi ^2/\nu `$ due to the two different background distributions. After all, it is a subjective judgment call on the part of the phenomenologist as to whether the goodness-of-fit is satisfactory. For large $`\nu `$, only when $`\chi ^2/\nu `$ starts approaching 1.5 does one really begin to start worrying about the model. For $`\nu 100`$, the error expected in $`\chi ^2/\nu `$ is $`0.14`$, so uncertainties in the renormalized $`\chi ^2/\nu `$ of the order of several percent play no critical role. The accuracy of the renormalized values is perfectly adequate for the purpose of judging whether to keep or discard a model.
In summary, extensive computer simulations for sifted data sets show that by combining the $`\chi ^2`$ parameter determinations with the corrected covariance matrix from the $`\chi ^2`$ fit, we obtain also a “robust” estimate of the errors, basically independent of both the background distribution and the model. Further, the renormalized $`\chi _{\mathrm{min}}^2/\nu `$ is a good predictor of the goodness-of-fit. Having to make a $`\mathrm{\Lambda }_0^2`$ fit to sift the data and then a $`\chi ^2`$ fit to the sifted data is a small computing cost to pay compared to the ability to make accurate predictions. Clearly, if the data are not badly contaminated with outliers, e.g., if a $`\mathrm{\Delta }\chi _i^2(x_i;𝜶)_{\mathrm{max}}=6`$ fit is satisfactory, the additional penalty paid is that the errors are enlarged by a factor of $`1.06`$ (see Table 1), which is not unreasonable to rescue a data set. Finally, if you are not happy about the error determinations, you can use the parameter estimates you have found to make Monte Carlo simulations of your model. By repeating a $`\mathrm{\Lambda }_0^2`$ fit to the simulated distributions and then sifting them to the same value of $`\mathrm{\Delta }\chi _i^2(x_i;𝜶)_{\mathrm{max}}`$ as was used in the initial determination of the parameters, and finally, by making a $`\chi ^2`$ fit to the simulated sifted set you can make an error determination based on the spread in the parameters found from the simulated data sets.
### 3.6 The parabola
As a final example of computer-generated data, we generated one noisy data set using a parabolic model. A total of 135 points were generated by computer. Using random number generators, the first 50 points generated picked $`x_i`$’s distributed randomly from 0 to 10. For each point $`x_i`$, a theoretical value $`\overline{y}_i`$ was found using $`\overline{y}_i=1+2x_i+0.5x_i^2`$. Next, the value of $`\sigma _i`$, the “experimental error”, i.e, the error bar assigned to point $`i`$, was generated randomly on the interval 0.2 to 2.7. The $`y_i`$’s were then generated, normally distributed about the value of $`\overline{y}_i`$ using the $`\sigma _i`$ that had been previously found. The next 50 points were chosen in the same manner, except that these $`\sigma _i`$ were randomly distributed between 0.2 and 5.2. This sample of 100 points made up the signal.
The 35 noise points were generated around a “nearby” parabola, given by $`\overline{y}_i=12+2x_i+0.2x_i^2`$. The first 15 points had their $`x_i`$ again randomly generated in the interval 0 to 10. The error bars assigned to each point were randomly distributed in the interval 0.2 to 5.2. To provide the outliers, the value of the theoretical $`\overline{y}_i`$ was found using a new parabola $`\overline{y}_i=12+2x_i+0.2x_i^2`$. These points were then normally distributed using $`\sigma _i`$’s uniformly distributed in the interval 0.8 to 20.8. The next 20 were generated in the same fashion, except that the error bars were uniformly distributed in the interval 0.2 to 8.2 and the $`y_i`$ values normally distributed with $`\sigma _i`$’s in the interval 1.6 to 65.6. In this case, we not only made “outliers”, but also contaminated the sample with substantial “inliers”, since we used a “nearby parabola” to generate the background data. Of course, this violates our Assumption 3 that we only have outliers, but gives us a feeling of what happens if substantial amounts of “inliers” are also present.
The resulting distribution of 135 points is shown in Fig. 8a, with the 100 squares being the normally distributed data and the 35 circles being the noise data.
The sifted data set, shown in Fig. 8b, retained 113 points after the $`\mathrm{\Delta }\chi _{i}^{2}{}_{\mathrm{max}}{}^{}=6`$ condition was applied to the original 135 points. At that point, we made both a conventional $`\chi ^2`$ fit to the sifted data set in order to evaluate the parameters, their errors and the goodness of fit. The $`\chi ^2`$ fit to the sifted data had $`\chi _{\mathrm{min}}^2=123.6`$, with $`\nu =110`$, giving $`\chi _{\mathrm{min}}^2/\nu =1.12`$. Renormalizing using $``$ found from eq. (5) , we get the corrected $`\times \chi _{\mathrm{min}}^2/\nu =1.24`$, whereas we expect $`1\pm 0.13`$. This is a reasonable fit with a probability of $`0.06`$. After using the Sieve algorithm, by minimizing $`\chi ^2`$, we found that the best-fit parabola, $`y=c_0+c_1x+c_2x^2`$, had $`c_0=1.18\pm 0.23`$ and $`c_1=2.05\pm 0.05`$ and $`c_2=0.489\pm 0.005`$, where the errors have been renormalized by the factor $`r_{\chi ^2}=1.05`$ found from eq. (4).
Figure 8b shows the results of using the Sieve procedure with the cut $`\mathrm{\Delta }\chi _i^2(x_i;𝜶)_{\mathrm{max}}=6`$. Of the original 135 points, all 100 of the signal points were retained (squares). There were 13 noise points (circles) also retained, all very close to the fitted straight line. These points are the “inliers” that resulted from the background generation using the “nearby parabola”, violating our primary assumption that there are only “outliers” as background. Thus, it is of great interest to see how well the Sieve procedure worked.
Had we applied a $`\chi ^2`$ minimization to original 130 point data set, we would have found $`\chi _{\mathrm{min}}^2/\nu =19.93`$, which clearly has infinitesimal statistical probability. The parabola resulting from this $`\chi ^2`$ fit is also shown in Fig. 8b. It clearly misses many of the data points in the sifted set.
When we fitted the parabola to only the 100 signal points, with no noise included, we got the parameters: $`c_0=0.97\pm 0.21,c_1=2.13\pm 0.05`$ and $`c_2=0.480\pm 0.005`$, using a conventional $`\chi ^2`$ fit. These parameters, within errors the same as those found using the “Sieve” algorithm, give a curve that is essentially indistinguishable from the solid line in Fig. 8b obtained using the Sieve algorithm. We note that even when the background produces some “inliers”, i.e., the cut $`\mathrm{\Delta }\chi _{i}^{2}{}_{\mathrm{max}}{}^{}`$ does not remove all of the background, the Sieve procedure is still very useful.
Finally, our procedure was completely stable for reasonable choices of $`\mathrm{\Delta }\chi _i^2`$, giving essentially the same answer for $`\mathrm{\Delta }\chi _i^2>4`$, 6 or 9. Thus, even in the presence of $`13\%`$ “inliers”, the answer after using the “Sieve” was reasonable. The parameter values are relatively unaffected, as are the errors. The main concern is the higher corrected $`\chi _{\mathrm{min}}^2/\nu `$ that is due to the background points that are close to the true signal and thus can not be “Sieved” out. However, this only affects the goodness-of-fit estimate, making $`\chi _{\mathrm{min}}^2/\nu `$ somewhat larger. In the end, the conclusion as to whether to accept the model or reject it on the basis of the goodness-of-fit estimate is a subjective judgment of the phenomenologist. Many models have been accepted when the $`\chi ^2`$ probability has been as low as a few tenths of a percent.
### 3.7 Real World data
We will illustrate the Sieve algorithm by simultaneously fitting all of the published experimental data above $`\sqrt{s}>6`$ GeV for both the total cross sections $`\sigma `$ and $`\rho `$ values for $`\overline{p}p`$ and $`pp`$ scattering, as well as for $`\pi ^{}p`$ and $`\pi ^+p`$ scattering. The $`\rho `$ value is the ratio of the real to the imaginary forward scattering amplitude and $`\sqrt{s}`$ is the cms energy $`E_{\mathrm{cms}}`$. The data sets used have been taken from the Web site of the Particle Data Group and have not been modified. They provide the energy ($`x_i`$), the measurement value ($`y_i`$) and the experimental error($`\sigma _i`$), assumed to be a standard deviation, for each experimental point.
Testing the hypothesis that the cross sections rise asymptotically as $`\mathrm{ln}^2s`$, as $`s\mathrm{}`$, the four functions $`\sigma ^\pm `$ and $`\rho ^\pm `$ that we will simultaneously fit for $`\sqrt{s}>6`$ GeV are:
$`\sigma ^\pm `$ $`=`$ $`c_0+c_1\mathrm{ln}\left({\displaystyle \frac{\nu }{m}}\right)+c_2\mathrm{ln}^2\left({\displaystyle \frac{\nu }{m}}\right)+\beta _𝒫^{}\left({\displaystyle \frac{\nu }{m}}\right)^{\mu 1}\pm \delta \left({\displaystyle \frac{\nu }{m}}\right)^{\alpha 1},`$ (7)
$`\rho ^\pm `$ $`=`$ $`{\displaystyle \frac{1}{\sigma ^\pm }}\left\{{\displaystyle \frac{\pi }{2}}c_1+c_2\pi \mathrm{ln}\left({\displaystyle \frac{\nu }{m}}\right)\beta _𝒫^{}\mathrm{cot}({\displaystyle \frac{\pi \mu }{2}})\left({\displaystyle \frac{\nu }{m}}\right)^{\mu 1}+{\displaystyle \frac{4\pi }{\nu }}f_+(0)\pm \delta \mathrm{tan}({\displaystyle \frac{\pi \alpha }{2}})\left({\displaystyle \frac{\nu }{m}}\right)^{\alpha 1}\right\},`$ (8)
$`{\displaystyle \frac{d\sigma ^\pm }{d(\nu /m)}}`$ $`=`$ $`c_1\left\{{\displaystyle \frac{1}{(\nu /m)}}\right\}+c_2\left\{{\displaystyle \frac{2\mathrm{ln}(\nu /m)}{(\nu /m)}}\right\}+\beta _𝒫^{}\left\{(\mu 1)(\nu /m)^{\mu 2}\right\}`$ (9)
$`\pm \delta \left\{(\alpha 1)(\nu /m)^{\alpha 2}\right\},`$
where the upper sign is for $`pp`$ ($`\pi ^+p`$) and the lower sign is for $`\overline{p}p`$ ($`\pi ^{}p`$) scattering. Here, $`\nu `$ is the laboratory energy of the projectile particle and $`m`$ is the proton (pion) mass. The exponents $`\mu `$ and $`\alpha `$ are real, as are the 6 constants $`c_0,c_1,c_2,\beta _𝒫^{},\delta `$ and the dispersion relation subtraction constant $`f_+(0)`$. We set $`\mu =0.5`$, appropriate for a Regge-descending trajectory, leaving us 7 parameters. We then require the fit to be anchored by the experimental values of $`\sigma _{\overline{p}p}`$ and $`\sigma _{pp}`$ ($`\sigma _{\pi ^{}p}`$ and $`\sigma _{\pi ^+p}`$), as well as their slopes, $`\frac{d\sigma ^\pm }{d(\nu /m)}`$, at $`\sqrt{s}=4`$ GeV for nucleon scattering and $`\sqrt{s}=2.6`$ GeV for pion scattering. This in turn imposes 4 conditions on the above equations and we thus have three free parameters to fit: $`c_1,c_2`$ and $`f_+(0)`$.
#### 3.7.1 $`\overline{p}p`$ and $`pp`$ scattering
The raw experimental data for $`\overline{p}p`$ and $`pp`$ scattering that are shown in Figures 1 and 2 were taken from the Particle Data Group. Figure 1 shows the $`\sigma _{\overline{p}p}`$ and $`\sigma _{pp}`$ data for $`\mathrm{E}_{\mathrm{cms}}>6`$ GeV, whereas Fig. 2 shows all of the experimental $`\rho _{\overline{p}p}`$ and $`\rho _{pp}`$ data for $`\mathrm{E}_{\mathrm{cms}}>6`$ GeV. There are a total of 218 points in these 4 data sets. We fit these 4 data sets simultaneously using eq. (7), eq. (8) and eq. (9). Before we applied the Sieve, we obtained $`\chi _{\mathrm{min}}^2=1185.6`$, whereas we expected 215. Clearly, either the model doesn’t work or there are a substantial number of outliers giving very large $`\mathrm{\Delta }\chi _i^2`$ contributions. The Sieve technique shows the latter to be the case.
We now study the effectiveness and stability of the Sieve. Table 2 contains the fitted results for $`\overline{p}p`$ and $`pp`$ scattering using 3 different choices of the cut-off, $`\mathrm{\Delta }\chi _{i}^{2}{}_{\mathrm{max}}{}^{}=4`$, 6 and 9. For each $`\mathrm{\Delta }\chi _{i}^{2}{}_{\mathrm{max}}{}^{}`$ cut it tabulates:
* the fitted parameters from the $`\chi ^2`$ fit together with the errors found in the $`\chi ^2`$ fit,
* the total $`\chi _{\mathrm{min}}^2`$,
* $`\nu `$, the number of degrees of freedom (d.f.) after the data have been sifted by the indicated $`\mathrm{\Delta }\chi _i^2`$ cut-off.
To get robust errors, the errors quoted in Table 2 for for each parameter should be multiplied by the common factor $`r_{\chi ^2}`$=1.05, from eq. (4), using the cut $`\mathrm{\Delta }=6`$.
We note that for $`\mathrm{\Delta }\chi _{i}^{2}{}_{\mathrm{max}}{}^{}=6`$, the number of retained data points is 193, whereas we started with 218, giving a background of $`13\%`$. We have rejected 25 outlier points (5 $`\sigma _{pp}`$, 5 $`\sigma _{\overline{p}p}`$, 15 $`\rho _{pp}`$ and no $`\rho _{\overline{p}p}`$ points) with $`\chi _{\mathrm{min}}^2`$ changing from 1185.6 to 182.8. We find $`\chi _{\mathrm{min}}^2/\nu =0.96`$, which when renormalized using eq. (5) for $`\mathrm{\Delta }=6`$ becomes $`\times \chi _{\mathrm{min}}^2/\nu =1.067`$, a very likely value with a probability of $`0.25`$.
Obviously, we have cleaned up the sample—we have rejected 25 datum points which had an average $`\mathrm{\Delta }\chi _i^240`$! We have demonstrated that: (1) the goodness-of-fit of the model is excellent, and (2) we had very large $`\mathrm{\Delta }\chi _i^2`$ contributions from the outliers that we were able to Sieve out. These outliers, in addition to giving a huge $`\chi _{\mathrm{min}}^2/\nu `$, severely distort the parameters found in a conventional $`\chi ^2`$ minimization, whereas they were easily handled by a robust fit which minimized $`\mathrm{\Lambda }_0^2`$, followed by a $`\chi ^2`$ fit to the sifted data.
Inspection of Table 2 shows that the parameter values $`c_1`$, $`c_2`$ and $`f_+(0)`$ effectively do not depend on $`\mathrm{\Delta }\chi _{i}^{2}{}_{\mathrm{max}}{}^{}`$, our cut-off choice, having only very small changes compared to the predicted parameter errors.
A further indication of the stability of the Sieve is illustrated in Table 3. As a function of $`\sqrt{s}`$, we have tabulated:
* the predicted total cross sections and $`\rho `$-values for $`\overline{p}p`$ and $`pp`$
* the errors in their predictions generated by the errors in the fit parameters $`c_1,c_2`$ and $`f_+(0)`$,
for two different cut-off values, $`\mathrm{\Delta }\chi _{i}^{2}{}_{\mathrm{max}}{}^{}=4`$ and 6. The predicted cross sections and $`\rho `$-values for the two values of $`\mathrm{\Delta }\chi _{i}^{2}{}_{\mathrm{max}}{}^{}`$ are virtually indistinguishable, giving us strong confidence in the Sieve technique when used with four different types of real-world experimental data.
The results of applying the Sieve algorithm to the 4 data sets, along with the fitted curves, are graphically shown in Fig. 10 for $`\sigma _{\overline{p}p}`$ and $`\sigma _{pp}`$ and in Fig. 11 for $`\rho _{\overline{p}p}`$ and $`\rho _{pp}`$. The total number of data points shown in Fig. 10 and in Fig. 11 is 193, whereas we started with 218 points. The fits shown are in excellent agreement with the 193 data points.
As a final test, we tried fitting another model which had its cross section energy dependence asymptotically rising as $`\mathrm{ln}s`$. This is the equivalent of setting the parameter $`c_2=0`$, leaving us two free parameters to fit, $`c_1`$ and $`f_+(0)`$. Using the same sifted data set which had given $`\chi _{\mathrm{min}}^2=182.8`$ for the $`\mathrm{ln}^2s`$ model we now obtained $`\chi _{\mathrm{min}}^2=1185.6`$ for only one more degree of freedom, clearly indicating that the $`\mathrm{ln}s`$ model was a very bad fit and could be excluded, whereas the $`\mathrm{ln}^2s`$ model gave a very good fit to the same data subset.
#### 3.7.2 $`\pi ^{}p`$ and $`\pi ^+p`$ scattering
The raw experimental data for $`\pi ^{}p`$ and $`\pi ^+p`$ scattering shown in Figures 3 and 4 were taken from the Particle Data Group. For $`\mathrm{E}_{\mathrm{cms}}>6`$ GeV, Figure 3 shows the $`\sigma _{\pi ^{}p}`$ and $`\sigma _{\pi ^+p}`$ data and Fig. 4 shows the $`\rho _{\pi ^{}p}`$ and $`\rho _{\pi +p}`$ data. There are a total of 155 points in these 4 data sets. Before we applied the Sieve algorithm, we obtained $`\chi ^2=527.8`$, whereas we expected 152, leading us to conclude that either the model doesn’t work or there are a substantial number of outliers giving very large $`\mathrm{\Delta }\chi _i^2`$ contributions. Once again, the Sieve technique shows the latter to be the case.
Table 4 contains the fitted results for $`\pi ^{}p`$ and $`\pi ^+p`$ scattering using 3 different choices of the cut-off, $`\mathrm{\Delta }\chi _{i}^{2}{}_{\mathrm{max}}{}^{}=4`$, 6 and 9. For each $`\mathrm{\Delta }\chi _{i}^{2}{}_{\mathrm{max}}{}^{}`$ it tabulates:
* the fitted parameters from the $`\chi ^2`$ fit together with the errors found in the $`\chi ^2`$ fit,
* the total $`\chi _{\mathrm{min}}^2`$,
* $`\nu `$, the number of degrees of freedom (d.f.) after the data have been sifted by the indicated $`\mathrm{\Delta }\chi _{i}^{2}{}_{\mathrm{max}}{}^{}`$ cut-off.
To get robust errors, the errors quoted in Table 4 for $`\mathrm{\Delta }\chi _i^2(x_i;𝜶)_{\mathrm{max}}=6`$ for each parameter should be multiplied by the common factor $`r_{\chi ^2}`$=1.05 of eq. (4) for the cut $`\mathrm{\Delta }=6`$.
For $`\mathrm{\Delta }\chi _{i}^{2}{}_{\mathrm{max}}{}^{}=6`$, the number of retained data points is 130, whereas we started with 155, a background of $`19\%`$. We have rejected 25 outlier points (2 $`\sigma _{\pi ^+p}`$, 19 $`\sigma _{\pi ^{}p}`$, 4 $`\rho _{\pi ^+p}`$ and no $`\rho _{\pi ^{}p}`$ points) with $`\chi _{\mathrm{min}}^2`$ changing from 527.8 to 148.1. We find $`\chi _{\mathrm{min}}^2/\nu =1.166`$, which when renormalized using eq. (5) for $`\mathrm{\Delta }=6`$ becomes $`\times \chi _{\mathrm{min}}^2/\nu =1.26`$, corresponding to a probability of 0.03, which is acceptable being about a $`2\sigma `$ effect.
Again, we have cleaned up the sample. We have rejected 25 datum points which had an average $`\mathrm{\Delta }\chi _i^215`$. We have demonstrated that: (1) the model works, and (2) we had large $`\mathrm{\Delta }\chi _i^2`$ contributions from the outliers that we were able to Sieve out.
Inspection of Table 4 shows that the parameter values effectively do not depend on our choice of cut-off, $`\mathrm{\Delta }\chi _{i}^{2}{}_{\mathrm{max}}{}^{}`$, not changing significantly compared to the predicted parameter errors. Another and perhaps better indication of the stability of the Sieve is illustrated in Table 5. Tabulated as a function of $`\sqrt{s}`$ are:
* the predicted total cross sections and $`\rho `$-values for $`\pi ^{}p`$ and $`\pi ^+p`$
* the errors in their predictions generated by the errors in the fit parameters $`c_1,c_2`$ and $`f_+(0)`$
for two different values of the cut-off, $`\mathrm{\Delta }\chi _{i}^{2}{}_{\mathrm{max}}{}^{}=4`$ and $`\mathrm{\Delta }\chi _{i}^{2}{}_{\mathrm{max}}{}^{}=6`$. The predicted cross sections and $`\rho `$ values for the two values of $`\mathrm{\Delta }\chi _{i}^{2}{}_{\mathrm{max}}{}^{}`$ are essentially indistinguishable, again generating strong confidence in the Sieve technique when used with these four different examples of real-world experimental data.
The results of applying the Sieve algorithm to the 4 data sets, along with the fitted curves, are graphically shown in Fig. 12 for $`\sigma _{\pi ^{}p}`$ and $`\sigma _{\pi ^+p}`$ and in Fig. 13 for $`\rho _{\pi ^{}p}`$ and $`\rho _{\pi ^+p}`$. The fits shown are in reasonable agreement with the 155 data points retained by the Sieve.
Again, when we attempted to fit the sifted data set of 130 points with a $`\mathrm{ln}s`$ fit, we found $`\chi _{\mathrm{min}}^2=942.5`$, with $`\nu =128`$, giving $`\chi ^2/\nu =7.35`$, with a probability of $`<<10^{45}`$. Thus, again a $`\mathrm{ln}^2s`$ fits well and a $`\mathrm{ln}s`$ fit is ruled out for the $`\pi p`$ system.
## 4 Comments and conclusions
We have shown that the Sieve algorithm works well in the case of backgrounds in the range of 0 to $`40\%`$, i.e., for extensive computer data that were generated about a straight line, as well as about a constant, and for a single event with a 20% outlier contamination as well as a 13%“inlier” contamination, that was generated about a parabola. It also works well for the $`13\%`$ to 19% contamination for the eight real-world data sets taken from the Particle Data Group. However, the Sieve algorithm is clearly inapplicable in the situation where the outliers (noise) swamps the signal. In that case, nothing can be done.
There are many possible choices for distributions resulting in robust fits. Our particular choice of minimizing the Lorentzian squared, $`\mathrm{\Lambda }_0^2(𝜶;𝒙)_{i=1}^N\mathrm{ln}\left\{1+0.18\mathrm{\Delta }\chi _i^2(x_i;𝜶)\right\}`$, in order to extract the parameters $`\{\alpha _1,\mathrm{},\alpha _M\}`$ needed to apply our Sieve technique seems to be a sensible one for both artificial computer-generated noisy distributions, as well as for real-world experimental data. This statement should not be interpreted as meaning that real-world data is truly well-approximated as a Lorentz distribution, but rather, as demonstrating that using the Lorentz distribution to get rid of outliers without sensibly affecting the fit parameters works well in the real world. Next, the choice of filtering out all points with $`\mathrm{\Delta }\chi _i^2>\mathrm{\Delta }\chi _{i}^{2}{}_{\mathrm{max}}{}^{}`$—where $`\mathrm{\Delta }\chi _{i}^{2}{}_{\mathrm{max}}{}^{}`$ is as large as possible—is optimal in both minimizing the loss of good data and maximizing the loss of outliers, resulting in a renormalized $`\times \chi _{\mathrm{min}}^2/\nu 1`$ for both the computer-generated and the real-world sample, as well as minimizing the distribution widths, and thus, the errors assigned to the parameters.
In detail, the utilization of the “Sieved” sample with $`\mathrm{\Delta }\chi _i^2<\mathrm{\Delta }\chi _{i}^{2}{}_{\mathrm{max}}{}^{}`$ allows one to
* use the unbiased parameter values found in a $`\chi ^2`$ fit to the truncated sample for the cut $`\mathrm{\Delta }\chi _i^2(x_i;𝜶)_{\mathrm{max}}`$, even in the presence of considerable background.
* find the renormalized $`\chi _{\mathrm{min}}^2/\nu `$, i.e., $`\times \chi _{\mathrm{min}}^2/\nu `$, where $``$ is the inverse of the factor given in eq. (5) as a function of $`\mathrm{\Delta }=\mathrm{\Delta }\chi _i^2(x_i;𝜶)_{\mathrm{max}}`$ and plotted in Figure 9.
* use the renormalized $`\chi _{\mathrm{min}}^2/\nu `$ to estimate the goodness-of-fit of the model employing the standard $`\chi ^2`$ probability distribution function. We thus estimate the probability that the data set fits the model, allowing one to decide whether to accept or reject the model.
* make a robust evaluation of the parameter errors and their correlations, by multiplying the standard covariance matrix $`C`$ found in the $`\chi ^2`$ fit by the appropriate value of $`(r_{\chi ^2})^2`$ for the cut $`\mathrm{\Delta }\chi _{i}^{2}{}_{\mathrm{max}}{}^{}`$. The value of $`r_{\chi ^2}`$ is given by eq. (4) and shown in Figure 9 as a function of the cut $`\mathrm{\Delta }\chi _{i}^{2}{}_{\mathrm{max}}{}^{}`$, where it is called $`\mathrm{\Delta }`$. It ranges from 1 for very large $`\mathrm{\Delta }`$ to $`1.14`$ for $`\mathrm{\Delta }=2`$ in eq. (4). However, this is not the complete story. The parameter error is $`\sigma =r_{\chi ^2}\times \mathrm{\Sigma }`$ and we must also take into account the increase in $`\mathrm{\Sigma }`$ due to the cut $`\mathrm{\Delta }`$, which causes the loss of signal points. As shown in Table 1 and discussed in detail in Section 3.4, the true loss of accuracy at $`\mathrm{\Delta }=2`$—relative to an unsifted sample of signal data—is the factor $`1.25`$. Thus, the algorithm starts failing rapidly for cuts $`\mathrm{\Delta }`$ smaller than 2.
In conclusion, the “ Sieve” algorithm gains its strength from the combination of making first a $`\mathrm{\Lambda }_0^2`$ fit to get rid of the outliers and then a $`\chi ^2`$ fit to the sifted data set. By varying the $`\mathrm{\Delta }\chi _i^2(x_i;𝜶)_{\mathrm{max}}`$ to suit the data set needs, we easily adapt to the different contaminations of outliers that can be present in real-world experimental data samples.
Not only do we now have a robust goodness-of-fit estimate, but we also have also a robust estimate of the parameters and, equally important, a robust estimate of their errors and correlations. The phenomenologist can now eliminate the use of possible personal bias and guesswork in “cleaning up” a large data set.
## 5 Acknowledgements
I would like to thank Professor Steven Block of Stanford University for valuable criticism and contributions to this manuscript and Professor Louis Lyons of Oxford University for many valuable discussions. Further, I would like to acknowledge the hospitality of the Aspen Center for Physics.
## Appendix A Robust Estimation
The terminology, “robust” statistical estimators, was first introduced to deal with small numbers of data points which have a large departure from the model predictions, i.e., outlier points. Later, research on robust estimation based on influence functions was carried out. More recently, robust estimations using regression models were made—these are inadequate for fitting non-linear models which often are needed in practical applications. For example, the fit needed for eq. (8) is a non-linear function of the coefficients $`c_0,c_1,c_2,\mathrm{}`$, since it is the ratio of two linear functions. We will discuss one possible technique for handling outlier points in a non-linear fit when we introduce the Lorentz probability density function in Section A.4.
### A.1 Maximum Likelihood Estimates
Let $`P_i`$ be the probability density of the $`i`$th individual measurement, $`i=1,\mathrm{},N`$, in the interval $`\mathrm{\Delta }y`$. Then the probability of the total data set is
$$𝒫=\underset{i=1}{\overset{N}{}}P_i\mathrm{\Delta }y.$$
(10)
Let us define the quantity
$$\mathrm{\Delta }\chi _i^2(x_i;𝜶)\left(\frac{y_iy(x_i;𝜶)}{\sigma _i}\right)^2,$$
(11)
where $`y_i`$ is the measured value at $`x_i`$, $`y(x_i;𝜶)`$ is the expected (theoretical) value from the model under consideration, and $`\sigma _i`$ is the experimental error of the $`i`$th measurement. The $`M`$ model parameters $`\alpha _k`$ are given by the $`M`$-dimensional vector $`𝜶=\{\alpha _1,\mathrm{},\alpha _M\}.`$
$`𝒫`$ is identified as the likelihood function, which we shall maximize as a function of the parameters $`𝜶=\{\alpha _1,\mathrm{},\alpha _M\}`$.
For the special case where the errors are normally distributed (Gaussian distribution), we have the likelihood function $`𝒫`$ given as
$$𝒫=\underset{i=1}{\overset{N}{}}\left\{\mathrm{exp}\left[\frac{1}{2}\left(\frac{y_iy(x_i;𝜶)}{\sigma _i}\right)^2\right]\frac{\mathrm{\Delta }y}{\sqrt{2\pi }\sigma _i}\right\}=\underset{i=1}{\overset{N}{}}\left\{\mathrm{exp}\left[\frac{1}{2}\mathrm{\Delta }\chi _i^2\right]\frac{\mathrm{\Delta }y}{\sqrt{2\pi }\sigma _i}\right\},$$
(12)
Maximizing the likelihood function $`𝒫`$ in eq. (12) is the same as minimizing the negative logarithm of $`𝒫`$, namely,
$$\underset{i=1}{\overset{N}{}}\frac{1}{2}\left(\frac{y_iy(xi;𝜶)}{\sigma _i}\right)^2N\mathrm{ln}\frac{\mathrm{\Delta }y}{\sqrt{2\pi }\sigma _i}.$$
(13)
Since $`N`$, $`\mathrm{\Delta }y`$ and $`\sigma _i`$ are constants, after using eq. (11), this is equivalent to minimizing the quantity
$$\frac{1}{2}\underset{i=1}{\overset{N}{}}\mathrm{\Delta }\chi _i^2(x_i;𝜶).$$
(14)
We now define $`\chi ^2(𝜶;𝒙)`$ as
$$\chi ^2(𝜶;𝒙)=\underset{i=1}{\overset{N}{}}\mathrm{\Delta }\chi _i^2(x_i;𝜶),$$
(15)
where $`𝒙\{x_1,\mathrm{},x_i,\mathrm{},x_N\}`$.
Hence, the $`\chi ^2`$ minimization problem, appropriate to the Gaussian distribution, reduces to
$$\mathrm{minimize}\mathrm{over}𝜶,\chi ^2(𝜶;𝒙)=\underset{i=1}{\overset{N}{}}\mathrm{\Delta }\chi _i^2(x_i;𝜶)$$
(16)
for the set of $`N`$ experimental points at $`x_i`$ having the value $`y_i`$ and error $`\sigma _i`$.
### A.2 Gaussian Distribution
To minimize $`\chi ^2`$, we must solve the (in general, non-linear) set of $`M`$ equations
$$\underset{i=1}{\overset{N}{}}\frac{1}{\sigma _i}\left(\frac{y_iy(x_i;𝜶)}{\sigma _i}\right)\left(\frac{y(x_i;\mathrm{}\alpha _j\mathrm{})}{\alpha _j}\right)=0,j=1,\mathrm{},M.$$
(17)
The Gaussian distribution allows a $`\chi ^2`$ minimization routine to return several exceedingly useful statistical quantities. Firstly, it returns the best-fit parameter space $`𝜶_{\mathrm{min}}`$. Secondly, the value of $`\chi _{\mathrm{min}}^2`$, when compared to the number of degrees of freedom ( d.f.$`\nu =NM`$, the number of data points minus the number of fitted parameters) allows one to make standard estimates of the goodness of the fit of the data set to the model used, using the $`\chi ^2`$ probability distribution function, given in standard texts, for $`\nu `$ degrees of freedom. Further, $`C^1`$, the $`M\times M`$ matrix of the partial derivatives at the minimum, given by
$$\left[C^1\right]_{jk}=\frac{1}{2}\left(\frac{^2\chi ^2}{\alpha _j\alpha _k}\right)_{𝜶=𝜶_{\mathrm{min}}},$$
(18)
allows us to compute the standard covariance matrix $`C`$ for the individual parameters $`\alpha _i`$, as well as the correlations between $`\alpha _j`$ and $`\alpha _k`$. Thus, when the errors are distributed normally, the $`\chi ^2`$ technique not only gives us the desired parameters $`𝜶_{\mathrm{min}}`$, but also furnishes us with statistically meaningful error estimates of the fitted parameters, along with goodness-of-fit information for the data to the chosen model—very valuable quantities for any model under consideration.
### A.3 Robust Distributions
We can generalize the maximum likelihood function of eq. (12), which is a function of the variable $`\frac{y_iy(x_i;𝜶)}{\sigma _i}`$, as
$$𝒫=\underset{i=1}{\overset{N}{}}\left\{\mathrm{exp}\left[\rho \left(\frac{y_iy(x_i;𝜶)}{\sigma _i}\right)\right]\mathrm{\Delta }y\right\},$$
(19)
where the function $`\rho \left(\frac{y_iy(x_i;𝜶)}{\sigma _i}\right)`$ is the negative logarithm of the probability density. Note that the statistical function $`\rho `$ used in this Appendix has nothing to do with the $`\rho `$-value used in eq. (8). Thus, we now have to minimize the generalization of eq. (14), i.e.,
$$\mathrm{minimize}\mathrm{over}𝜶,\underset{i=1}{\overset{N}{}}\rho \left(\frac{y_iy(x_i;𝜶)}{\sigma _i}\right),$$
(20)
for the $`N`$-dimensional vector $`𝒙`$.
This yields the more general set of $`M`$ equations
$$\underset{i=1}{\overset{N}{}}\frac{1}{\sigma _i}\psi \left(\frac{y_iy(x_i;𝜶)}{\sigma _i}\right)\left(\frac{y(x_i;\mathrm{}\alpha _j\mathrm{})}{\alpha _j}\right)=0,j=1,\mathrm{},M,$$
(21)
where the influence function $`\psi (z)`$ in eq. (21) is given by
$$\psi (z)\frac{d\beta (z)}{dz},z\frac{y_iy(x_i;𝜶)}{\sigma _i}=\mathrm{sign}(y_iy(x_i;𝜶))\times \sqrt{\mathrm{\Delta }\chi _i^2(x_i;𝜶)}.$$
(22)
Comparison of eq. (21) with the Gaussian equivalent of eq. (17) shows that
$$\rho (z)=\frac{1}{2}z^2,\psi (z)=z(\mathrm{for}\mathrm{a}\mathrm{Gaussian}\mathrm{distribution}).$$
(23)
We note that for a Gaussian distribution, the influence function $`w(z)`$ for each experimental point $`i`$ is proportional to $`\sqrt{\mathrm{\Delta }\chi _i^2}`$, the normalized departure of the point from the theoretical value. Thus, the more the departure from the theoretical value, the more “influence” the point has in minimizing $`\chi ^2`$. This gives outliers (points with large departures from their theoretical values) unduly large “influence” in computing the best vector $`𝜶`$, easily skewing the answer due to the inclusion of these outliers.
### A.4 Lorentz Distribution
Consider the normalized Lorentz probability density distribution (also known as the Cauchy distribution or the Breit-Wigner line width distribution), given by
$`P(z)`$ $`=`$ $`{\displaystyle \frac{\sqrt{\gamma }}{\pi }}{\displaystyle \frac{1}{1+\gamma z^2}},`$ (24)
where $`\gamma `$ is a constant whose significance will be discussed later. Using eq. (11) and eq. (22), we rewrite eq. (24) in terms of the measurement errors $`\sigma _i`$ and the experimental measurements $`y_i`$ at $`x_i`$ as
$`P\left({\displaystyle \frac{y_iy(x_i;𝜶)}{\sigma _i}}\right)`$ $`=`$ $`{\displaystyle \frac{\sqrt{\gamma }}{\pi }}{\displaystyle \frac{1}{1+\gamma \left(\frac{y_iy(x_i;𝜶)}{\sigma _i}\right)^2}}`$ (25)
$`=`$ $`{\displaystyle \frac{\sqrt{\gamma }}{\pi }}{\displaystyle \frac{1}{1+\gamma \mathrm{\Delta }\chi _i^2(x_i;𝜶)}}.`$
It has long tails and therefore is more suitable for robust fits than is the Gaussian distribution. Taking the negative logarithm of eq. (25) and using it in eq. (20), we see that
$`\rho (z)`$ $`=`$ $`\mathrm{ln}\left(1+\gamma z^2\right)=\mathrm{ln}\left\{1+\gamma \mathrm{\Delta }\chi _i^2(x_i;𝜶)\right\}\mathrm{and}`$
$`\psi (z)`$ $`=`$ $`{\displaystyle \frac{z}{1+\gamma z^2}}={\displaystyle \frac{\mathrm{sign}(y_iy(x_i;𝜶))\times \sqrt{\mathrm{\Delta }\chi _i^2(x_i;𝜶)}}{1+\gamma \mathrm{\Delta }\chi _i^2(x_i;𝜶)}}.`$ (26)
In analogy to $`\chi ^2`$ minimization, we must now minimize $`\mathrm{\Lambda }^2(𝜶;𝒙)`$, the Lorentzian squared, with respect to the parameters $`𝜶`$, for a given set of experimental points $`𝒙`$, i.e.,
$$\mathrm{minimize}\mathrm{over}𝜶,\mathrm{\Lambda }^2(𝜶;𝒙)\underset{i=1}{\overset{N}{}}\mathrm{ln}\left\{1+\gamma \mathrm{\Delta }\chi _i^2(x_i;𝜶)\right\},$$
(27)
for the set of $`N`$ experimental points at $`x_i`$ having the value $`y_i`$ and error $`\sigma _i`$.
We have made extensive computer simulations using Gaussianly generated data (constant and straight line models) which showed empirically that the choice $`\gamma =0.18`$ minimized the rms (root mean square) parameter widths found in $`\mathrm{\Lambda }^2`$ minimization. Further, it gave rms widths that were almost as narrow as those found in $`\chi ^2`$ minimization on the same data. We will adopt this value of $`\gamma `$, since it effectively minimizes the width for the $`\mathrm{\Lambda }^2`$ routine, which we now call $`\mathrm{\Lambda }_0^2(𝜶;𝒙)`$. Thus we select for our robust algorithm,
$$\mathrm{minimize}\mathrm{over}𝜶,\mathrm{\Lambda }_0^2(𝜶;𝒙)\underset{i=1}{\overset{N}{}}\mathrm{ln}\left\{1+0.18\mathrm{\Delta }\chi _i^2(x_i;𝜶)\right\}.$$
(28)
An important property of $`\mathrm{\Lambda }_0^2(𝜶;𝒙)`$ is that it numerically gives the same total $`\chi _{0_{\mathrm{min}}}^2`$ as that found in a $`\chi ^2`$ fit, i.e. $`\chi _0^2=_{i=1}^N\mathrm{\Delta }\chi _i^2(x_i;𝜶)`$, where the $`\mathrm{\Delta }\chi _i^2(x_i;𝜶)`$ come from the minimization of $`\mathrm{\Lambda }_0^2`$ in eq. (28), is the same as the $`\chi _{\mathrm{min}}^2`$ found using a standard $`\chi ^2`$ minimization on the same data.
We note from eq. (26) that the influence function for a point $`i`$ for small $`\sqrt{\mathrm{\Delta }\chi _i^2}`$ increases proportional to $`\sqrt{\mathrm{\Delta }\chi _i^2}`$ (just like the Gaussian distribution does), whereas for large $`\sqrt{\mathrm{\Delta }\chi _i^2}`$, it decreases as $`1/\sqrt{\mathrm{\Delta }\chi _i^2}`$. Thus, large outliers have much less “influence” on the fit than do points close to the model curve—this feature makes $`\mathrm{\Lambda }^2`$ minimization robust. Thus, outliers have little influence on the choice of the parameters $`𝜶_{\mathrm{min}}`$ resulting from the minimization of $`\mathrm{\Lambda }_0^2`$, a major consideration for a robust minimization method.
Unlike the minimization of $`\chi ^2`$, the minimization of $`\mathrm{\Lambda }_0^2`$, while yielding the desired robust estimate of $`𝜶_{\mathrm{min}}`$, gives neither parameter error information on $`𝜶_{\mathrm{min}}`$ nor a conventional goodness-of-fit. These are major failings, since one has no objective grounds for accepting or rejecting the model. We will rectify these shortcomings in the main section of the text, Section 2, where we describe the adaptive “Sieve” algorithm. Extensive computer studies, summarized in Section 3.4, demonstrate that use of this algorithm enables one to make a robust error estimate of $`𝜶_{\mathrm{min}}`$, as well as a robust estimate of the goodness-of-fit of the data to the model. |
warning/0506/nucl-th0506004.html | ar5iv | text | # The Ramsauer model for the total cross sections of neutron nucleus scattering
## Abstract
Theoretical study of systematics of neutron scattering cross sections on various materials for neutron energies up to several hundred MeV are of practical importance. In this paper, we analysed the experimental neutron scattering total cross sections from 20MeV to 550MeV using Ramsauer model for nuclei ranging from Be to Pb.
total cross section, neutron nucleus scattering, Ramsauer model, empirical formulae.
In recent times world over, there is a renewed interest in the neutron nucleus scattering data. This is owing to a new concept of controlled nuclear energy source called the accelerator driven sub-critical (ADS) system ads1 ; ads2 . In this ADS system, neutrons are produced by bombarding a heavy element target with a high energy proton beam of typically above 1.0GeV with a current of $`>10mA`$ ads1 . The sub-critical reactor is driven critical by spallation neutrons produced by proton beam on typically a molten lead target. Such accelerator driven systems can also be used for waste incineration of the long lived radio active waste produced in reactors based on thermal neutron induced fission. Reactor physics calculations of these ADS type of systems require neutron-nucleus scattering cross sections up to 500 MeV neutron energy. Unlike the neutron energy spectrum from a thermal neutron fission, the spallation neutron energies reach up to several hundred MeV. Therefore, it is currently very important to study the systematics of neutron scattering cross sections on various nuclei for neutron energies up to several hundred MeV. In the present work, we performed an empirical study of the energy dependence of total cross sections ($`\sigma _{tot}`$) of the neutron-nucleus (n-N) scattering. It is well known that the total cross sections are explained by the nuclear Ramsauer model. The Ramsauer model was first proposed by Lawson lawson in the year 1953 to phenomenologically explain the energy dependence of total cross sections of neutron nucleus scattering. In order to have insight of the working of this model, it is necessary to understand the optical model (OM) description of neutron scattering. In the OM approach, complex optical model potentials (OMP) are used and the Schrodinger’s equation is solved to obtain the scattering amplitude. The real part of the OMP describes the scattering and the imaginary part results in attenuation or absorption of the incident wave. The reaction cross section is given by the absorption of the neutron flux. The scattering calculations are performed using partial wave expansion method and the phase shifts ($`\eta _{\mathrm{}}=\alpha _{\mathrm{}}e^{i\beta _{\mathrm{}}})`$ are determined. These complex phase sifts are strongly angular momentum and energy dependent for a given set of potentials. In terms of the phase shifts and the wave number ($`\mathrm{\lambda ̄}=\mathrm{}/\sqrt{2mE}`$), the total cross sections are given below.
$`\sigma _{tot}`$ $`=`$ $`2\pi \mathrm{\lambda ̄}^2{\displaystyle \underset{\mathrm{}}{}}(2\mathrm{}+1)\left[1\mathrm{}\eta _{\mathrm{}}\right]`$ (1)
Extensive study of the optical model fits of scattering cross sections on various nuclei over wide energy range have been made by several groups. This is owing to the excellent data base of neutron total cross sections available in the energy range up to 600 MeV finlay ; dietrich1 ; abfal . The most recent work by Koning and Delaroche (KD) kd presents a very exhaustive search for OMP parameters that fit the data very well up to 200 MeV. Alternatively, the nuclear Ramsauer model lawson provides a simple means to parameterise the energy dependence of neutron nucleus total scattering cross sections. This model assumes that the scattering phase shifts are independent of angular momentum ($`\mathrm{}`$) as given in Eq.(2) ($`\eta =\alpha e^{i\beta })`$ , in contrast to the predictions of the optical model given in Eq.(1). Further, it was proposed that the $`\mathrm{}`$-independent phase shift varies slowly with energy. This model was successfully applied for neutron scattering from various nuclei by Peterson peterson ; book . There were some attempts franco ; gould ; anderson ; grimes1 (see references therein) to put this Ramsauer model on a sound theoretical basis. The neutron total cross sections have thus been well studied using this model, over a wide range of nuclear masses as well as neutron energies up to 500 MeV anderson ; bauer ; madsen ; grimes1 ; grimes2 ; grimes3 ; dietrich2 . Deb et. al., deb have achieved simple functional forms for the total cross sections by parameterising the maximum partial wave ($`\mathrm{}_0=\mathrm{}_{max}`$) values.
In our earlier work surya1 , we presented the Ramsauer model analysis of the results of optical model code SCAT2 scat2 using KD potentials. In the present work, we performed the Ramsauer model analysis of the experimental data of neutron total cross sections for heavy and light nuclei by using Eq.(2). The quantities $`R(fm),\alpha ,\beta `$ are functions of atomic mass number (A) and the center of mass energy (E).
$`\sigma _{tot}`$ $`=`$ $`2\pi (R+\mathrm{\lambda ̄})^2\left(1\alpha \mathrm{cos}\beta \right)`$ (2)
$`\beta `$ $`=`$ $`\beta _xA^{\frac{1}{3}}(\sqrt{E+V}\sqrt{E})`$ (3)
$`V`$ $`=`$ $`V_A+V_E\sqrt{E}`$
$`V_A`$ $`=`$ $`V_0+V_1(NZ)/A+V_2/A`$
$`\alpha `$ $`=`$ $`\alpha _0+\alpha _A\sqrt{E}`$
$`\alpha _A`$ $`=`$ $`\alpha _1\mathrm{ln}(A)+\alpha _2/\mathrm{ln}(A)`$
$`R`$ $`=`$ $`r_0A^{\frac{1}{3}}+r_A\sqrt{E}+r_2`$
$`r_A`$ $`=`$ $`r_{10}\mathrm{ln}(A)+r_{11}/\mathrm{ln}(A)`$
$`r_0`$ $`=`$ $`1.42988,r_{10}=0.02298,r_{11}=0.10268`$ (4)
$`r_2`$ $`=`$ $`0.23216,V_0=46.51099,V_1=6.73833`$
$`V_2`$ $`=`$ $`117.52082,V_E=3.21817,\beta _x=0.5928`$
$`\alpha _0`$ $`=`$ $`0.02868,\alpha _1=0.00274,\alpha _2=0.13211`$
Figure 1 shows the Ramsauer model fits (solid lines) for $`\sigma _{tot}`$ cross sections using Eqs.(2),(3),(4) and the symbols represent the experimental data. The fits are obtained with total of twelve free parameters as given in Eq.(4), over wide mass range of $`{}_{}{}^{24}Mg`$ to $`{}_{}{}^{208}Pb`$. In Fig.2, we show the Ramsauer model fits for light nuclei from <sup>9</sup>Be to <sup>24</sup>Mg. These fits cover the neutron energy region ($`E_{cm}`$) of 20-550 MeV. Similar Ramsauer model fits to total cross sections were already shown by various groups grimes1 ; grimes2 ; deb (see the references therein). As shown in Figs. (1,2), the functional dependence on energy and mass given in Eqs.(2),(3) with twelve global parameters was able to reproduce the experimental data very well.
As shown in Fig. 2, the parameters are same as those for heavy systems, except for the radius parameter. The radius parameter for scattering from light nuclei turns out to be rather larger than for heavy nuclei. In some cases of neutron nucleus scattering, the radius parameter (only r<sub>0</sub>) has to be specially adjusted to achieve best fits. These cases did not obey the Ramsauer model systematics for either light or heavy nuclei. In case of <sup>27</sup>Al, <sup>232</sup>Th and <sup>238</sup>U, the adjusted radius parameters (in fm) are respectively, r<sub>0</sub>= 1.3892(Al), 1.4488(Th), 1.44641(U). These values are comparable to the general systematics of r<sub>0</sub> values for heavy nuclei (1.42988fm) and for light nuclei (1.53091fm).
In conclusion, we performed the Ramsauer model parameterization of experimental neutron total scattering cross sections from light and heavy nuclei. The parameters have been found to be same for light and heavy nuclei, except for the radius parameters. We could reproduce the experimental neutron total cross sections by using twelve parameters in the energy range of 20MeV to 550 MeV. We proposed a new functional form for energy dependence and atomic mass dependence of the Ramsauer model parameters. |
warning/0506/math0506534.html | ar5iv | text | # On deformations of Q-factorial symplectic varieties
## 1 Introduction
This is a supplement to \[Na\]. In \[Na\] we have proposed a new category of complex symplectic varieties which admit certain singularities. It includes symplectic V-manifolds (cf. \[Fu 1\]) and O’Grady’s singular moduli spaces of semi-stable torsion free sheaves on a K3 surface \[O\] (cf. \[Na\], Introduction). A normal compact Kähler space $`Z`$ is a symplectic variety if its regular locus $`U`$ admits a non-degenerate holomorphic closed 2-form $`\omega `$ with the following property: for any resolution $`\nu :\stackrel{~}{Z}Z`$ of $`Z`$ with $`\nu ^1(U)U`$, the 2-form $`\omega `$ on $`\nu ^1(U)`$ extends to a holomorphic 2-form on $`\stackrel{~}{Z}`$. If $`Z`$ can be embedded in a projective space, $`Z`$ is called a projective symplectic variety.
Our purpose is to give a positive answer to the following problem posed in \[Na\]:
Problem. Let $`Z`$ be a Q-factorial projective symplectic variety with terminal singularities. Assume that $`Z`$ is smoothable by a suitable flat deformation. Is $`Z`$ then non-singular from the first ?
We actually prove more:
Main Theorem. Let $`Z`$ be a Q-factorial projective symplectic variety with terminal singularities. Then any flat deformation of $`Z`$ is locally trivial; in other words, it preserves all singularities on $`Z`$.
This theorem would explain why “local Torelli theorem” \[Na, Theorem 8\] holds for such symplectic varieties.
Example (\[O\], cf. also \[K-L-S\]): Let $`S`$ be a K3 surface with Picard number $`1`$. Fix an even integer $`c`$ with $`c6`$. Let $`M_c`$ be the moduli space of semi-stable sheaves on $`S`$ of rank 2 with $`c_1=0`$ and $`c_2=c`$. O’Grady proved in \[O\] that $`M_c`$ is a projective symplectic variety with terminal singularities. In the preprint version of \[O\]: math.AG/9708009, he constructed a projective resolution $`\widehat{\pi }:\widehat{M}_cM_c`$. Let $`N_1(\widehat{M}_c/M_c)`$ be the R-vector space of $`\widehat{\pi }`$-numerical classes of 1-cycles $`l`$ on $`\widehat{M}_c`$ such that $`\widehat{\pi }(l)=\mathrm{point}`$. Then
$$N_1(\widehat{M}_c/M_c)=𝐑[\widehat{ϵ}_c]𝐑[\widehat{\sigma }_c]𝐑[\widehat{\gamma }_c]$$
for effective 1-cycles $`\widehat{ϵ}_c`$, $`\widehat{\sigma }_c`$ and $`\widehat{\gamma }_c`$ (ibid. Section 3). On the other hand, $`\mathrm{Exc}(\widehat{\pi })`$ consists of three prime divisors $`\widehat{\mathrm{\Omega }}_c`$, $`\widehat{\mathrm{\Sigma }}_c`$ and $`\widehat{\mathrm{\Delta }}_c`$ (ibid. Section 3). This implies that $`M_c`$ is Q-factorial. In fact, let $`D`$ be a prime divisor of $`M_c`$ and let $`D^{}`$ be its proper transform on $`\widehat{M}_c`$. By the observation above, we can find $`a_i𝐐`$ $`(i=1,2,3)`$ such that $`D^{\prime \prime }:=D^{}+a_1\widehat{\mathrm{\Omega }}_c+a_2\widehat{\mathrm{\Sigma }}_c+a_3\widehat{\mathrm{\Delta }}_c`$ is $`\widehat{\pi }`$-numerically trivial. Since $`M_c`$ has only rational singularities, there is a positive integer $`m`$ such that $`mD^{\prime \prime }`$ is linearly equivalent to the pull-back of a Cartier divisor $`E`$ of $`M_c`$. Then $`mD`$ is linearly equivalent to $`E`$. More generally, for a K3 surface or an Abelian surface, almost all moduli spaces of semi-stable sheaves with non-primitive Mukai vectors are Q-factorial projective symplectic varieties \[K-L-S\]. Applying Main Theorem to these moduli spaces, we conclude that any deformation of them are locally rigid; in particular, they have no smoothing via deformations.
In the final section, we shall apply Main Theorem to get the following remarkable results.
Corollary 1. Let $`h:YZ`$ and $`h^{}:Y^{}Z`$ be two Q-factorial terminalizations of a projective symplectic variety $`Z`$; in other words, $`h`$ (resp. $`h^{}`$) is a crepant projective birational morphism such that $`Y`$ (resp. $`Y^{}`$) has only Q-factorial terminal singularities. Assume that $`Y`$ is non-singular. Then $`Y^{}`$ is also non-singular.
Remark 1. A Q-factorial terminalization $`Y`$ of a projective symplectic variety $`Z`$ is again symplectic. In fact, there is a 2-form $`h^{}\omega `$ on $`h^1(U)`$, where $`U:=Z_{reg}`$. By the definition of a symplectic variety, this 2-form extends to a holomorphic 2-form on $`V:=Y_{reg}`$. Moreover, since $`K_Y=h^{}K_Z`$, the extended 2-form is non-degenerate over $`V`$. If we take a resolution $`\stackrel{~}{Y}`$ of $`Y`$, then this 2-form further extends to a holomorphic 2-form on $`\stackrel{~}{Y}`$
Example (flops): When the codimensions of $`\mathrm{Exc}(h)`$ and $`\mathrm{Exc}(h^{})`$ are both larger than one, Corollary 1 says that the smoothness is preserved under a flop. A similar result has been proved by Kaledin under certain conditions when $`Z`$ is an affine symplectic singularity (cf. “Symplectic resolutions: deformations snd birational maps”, math.AG/0012008).
Corollary 2. Let $`Z`$ be a projective symplectic variety. Assume that $`Z`$ has a Q-factorial terminalization $`h:YZ`$. Then the following are equivalent.
(1) $`Z`$ is smoothable by a flat deformation;
(2) $`Y`$ is non-singular.
The following conjecture would be true if any projective symplectic variety had a Q-factorial terminalization (cf. Remark 2, (1))
Conjecture. A projective symplectic variety is smoothable by a deformation if and only if it admits a crepant (symplectic) resolution.
Remark 2. (1) The existence of a Q-factorial terminalization follows from the “minimal model conjecture”. In fact, let $`\stackrel{~}{Z}Z`$ be a projective resolution. If the minimal model conjecture holds, one can get a Q-factorial terminalization $`YZ`$ after running the minimal model program for $`\stackrel{~}{Z}`$ over the fixed base $`Z`$. The Q-factorial terminalizations are not unique; however, they have the same kind of singularities by virtue of the proof of Corollary 1. The minimal model conjecture is proved by Shokurov \[Sh\] in dimension 4.
(2) It would be interesting to know whether our results also hold in the case where $`Z`$ is a convex symplectic variety. Here we say that $`Z`$ is convex if there is a birational projective morphism from $`Z`$ to an affine symplectic singularity $`X`$. When $`X`$ has non-isolated singularity, the usual deformation theory of $`Z`$ (and $`X`$) has an infinite dimensional nature and does not work well. Instead, we should consider a “Poisson deformation” (cf. \[G-K\]).
## 2 Proof of Main Theorem
(i) Q-factoriality: A normal variety $`Z`$ is called Q-factorial if, for any Weil divisor $`D`$ of $`Z`$, $`mD`$ is a Cartier divisor of $`Z`$ for some positive integer $`m`$. Assume that $`Z`$ is a projective variety with rational singularities. Let $`\nu :\stackrel{~}{Z}Z`$ be a resolution such that the exceptional locus consists of finite number of divisors $`E_i`$ $`(1im)`$. Denote by $`[E_i]H^2(\stackrel{~}{Z},𝐙)`$ the corresponding classes. By \[Ko-Mo, (12.1.6)\], $`Z`$ is Q-factorial if and only if
$$\mathrm{im}[H^2(\stackrel{~}{Z},𝐐)H^0(Z,R^2\nu _{}𝐐)]=\mathrm{im}[𝐐[E_i]H^0(Z,R^2\nu _{}𝐐)].$$
We give here a proof of the “only if” part. Fix an ample class $`\kappa H^2(\stackrel{~}{Z},𝐐)`$. Let $`H_0^2(\stackrel{~}{Z},𝐐)`$ be the primitive part of $`H^2(\stackrel{~}{Z},𝐐)`$. By definition, we have
$$H^2(\stackrel{~}{Z},𝐐)=H_0^2(\stackrel{~}{Z},𝐐)𝐐\kappa .$$
Since $`Z`$ has only rational singularities, the natural map $`H^2(Z,𝐐)H^2(\stackrel{~}{Z},𝐐)`$ is an injection of mixed Hodge structures (\[De\]); hence, by the strict compatibility of weight filtrations, $`H^2(Z,𝐐)`$ admits a pure Q-Hodge structure of weight 2. Then $`H_0^2(Z,𝐐):=H^2(Z,𝐐)H_0^2(\stackrel{~}{Z},𝐐)`$ also admits a pure Q-Hodge structure of weight 2. Note that $`H_0^2(Z,𝐂)`$ contains $`H^{2,0}(\stackrel{~}{Z})`$ and $`H^{0,2}(\stackrel{~}{Z})`$. With respect to the Hodge-Riemann bilinear form on $`H_0^2(\stackrel{~}{Z},𝐐)`$, we put $`V:=(H_0^2(Z,𝐐))^{}`$. One can check that
$$VH_0^2(Z,𝐐)=0,$$
and hence
$$H^2(\stackrel{~}{Z},𝐐)=V𝐐\kappa H_0^2(Z,𝐐).$$
Note that $`H_0^2(Z,𝐐)`$ is mapped to zero by the map $`H^2(\stackrel{~}{Z},𝐐)H^0(Z,R^2\nu _{}𝐐)`$. Since every element of $`V`$ is of type (1,1), we see that
$$\mathrm{im}[H^2(\stackrel{~}{Z},𝐐)H^0(Z,R^2\nu _{}𝐐)]=\mathrm{im}[\mathrm{Pic}(\stackrel{~}{Z})𝐐H^0(Z,R^2\nu _{}𝐐)].$$
Now take a Q-Cartier divisor $`D`$ of $`\stackrel{~}{Z}`$. The push-forward $`\nu _{}D`$ is a Q-Cartier divisor of $`Z`$ since $`Z`$ is Q-factorial. Now the pull-back $`\nu ^{}(\nu _{}D)`$ of $`\nu _{}D`$ is well-defined, and one has
$$D=\nu ^{}(\nu _{}D)+\mathrm{\Sigma }a_iE_i$$
for some $`a_i𝐐`$. This implies the “only if” part.
(ii) Terminal singularity: We recall the definition of a terminal singularity. Assume, for simplicity, that $`Z`$ is a normal variety such that the canonical divisor $`K_Z`$ is a Cartier divisor. We say that $`Z`$ has only terminal (resp. canonical) singularities if, for a resolution $`\nu :\stackrel{~}{Z}Z`$ such that the exceptional locus consists of finite number of divisors $`E_i`$ $`(1im)`$,
$$K_{\stackrel{~}{Z}}=\nu ^{}K_Z+\mathrm{\Sigma }a_iE_i$$
with $`a_i>0`$ (resp. $`a_i0`$) for all $`i`$. Let $`\mathrm{\Sigma }Z`$ be the singular locus. If $`Z`$ has only terminal singularities, then we have $`\mathrm{codim}(\mathrm{\Sigma }Z)3`$ (cf. \[Re\]). But if, in addition, $`Z`$ is a symplectic variety, then we have $`\mathrm{codim}(\mathrm{\Sigma }Z)4`$ by \[Na 2\], \[K\]. Conversely, if $`Z`$ is a symplectic variety with $`\mathrm{codim}(\mathrm{\Sigma }Z)4`$, then $`Z`$ has only terminal singularities \[Na 2\].
(iii) Infinitesimal deformations: Assume that $`Z`$ is a reduced complex space. Recall that a deformation $`𝒵`$ of $`Z`$ over an analytic space $`S`$ with a reference point $`0S`$ is a flat map $`𝒵S`$ with a fixed isomorphism $`𝒵\times _S\{0\}Z`$. When $`0S`$ is a spectrum of an Artinian local C-algebra with a unique point, $`𝒵`$ is called an infinitesimal deformation of $`Z`$. Moreover, an infinitesimal deformation $`𝒵S`$ is called locally trivial if, for any Stein open set $`VZ`$, $`𝒵|_VV\times S`$.
For $`n0`$, we put $`A_n:=𝐂[t]/(t^{n+1})`$ and $`S_n:=\mathrm{Spec}A_n`$. Suppose that $`Z_nS_n`$ is an infinitesimal deformation of $`Z`$ over $`S_n`$. Define $`T_{Z_n/S_n}:=\underset{¯}{\mathrm{Hom}}(\mathrm{\Omega }_{Z_n/S_n}^1,𝒪_{Z_n}).`$ Let $`S_n[ϵ]:=\mathrm{Spec}(A_n[ϵ]/ϵ^2)`$. There is a natural closed immersion $`S_nS_n[ϵ]`$. Moreover, the injection $`A_nA_n[ϵ]/ϵ^2`$ induces a map $`S_n[ϵ]S_n`$. $`\mathrm{Ext}^1(\mathrm{\Omega }_{Z_n/S_n}^1,𝒪_{Z_n})`$ plays an important role in the deformation theory of a reduced complex space:
(iii-a): There is a one-to-one correspondence between $`\mathrm{Ext}^1(\mathrm{\Omega }_{Z_n/S_n}^1,𝒪_{Z_n})`$ and the set of equivalence classes of infinitesimal deformations of $`Z_n`$ over $`S_n[ϵ]`$.
Here two infinitesimal deformations $`𝒵_n`$ and $`𝒵_n^{}`$ of $`Z_n`$ over $`S_n[ϵ]`$ are equivalent if there is an $`S_n[ϵ]`$-isomorphism between them which induces, over $`S_n`$, the identity map of $`Z_n`$.
The correspondence is given in the following manner. Assume an infinitesimal deformation $`𝒵_n`$ is given. By the map $`S_n[ϵ]S_n`$, we regard $`𝒵_n`$ as an $`S_n`$-analytic space. There is an exact sequence
$$ϵ𝒪_{Z_n}\stackrel{d}{}\mathrm{\Omega }_{𝒵_n/S_n}^1_{A_n[ϵ]}A_n\mathrm{\Omega }_{Z_n/S_n}^10.$$
Here $`d`$ is defined by $`d(ϵf):=fdϵ`$ for $`ϵfϵ𝒪_{Z_n}`$. By using the fact that $`Z`$ is reduced, one can show that this $`d`$ is an injection; hence we have an element of $`\mathrm{Ext}^1(\mathrm{\Omega }_{Z_n/S_n}^1,ϵ𝒪_{Z_n})`$. Since $`𝒵_n`$ is flat over $`S_n[ϵ]`$, $`ϵ𝒪_{Z_n}𝒪_{Z_n}`$.
Conversely, let there be an element of $`\mathrm{Ext}^1(\mathrm{\Omega }_{Z_n/S_n}^1,𝒪_{Z_n})\mathrm{Ext}^1(\mathrm{\Omega }_{Z_n/S_n}^1,ϵ𝒪_{Z_n})`$. This element gives an $`𝒪_{Z_n}`$-module extension of $`\mathrm{\Omega }_{Z_n/S_n}^1`$ by $`ϵ𝒪_{Z_n}`$. We consider this extension as an $`A_n`$-module extension. Then by the homomorphism $`\mathrm{Ext}_{A_n}^1(\mathrm{\Omega }_{Z_n/S_n}^1,ϵ𝒪_{Z_n})\mathrm{Ext}_{A_n}^1(𝒪_{Z_n},ϵ𝒪_{Z_n})`$ induced by $`d:ϵ𝒪_{Z_n}\mathrm{\Omega }_{Z_n/S_n}^1`$, we have the commutative diagram of exact sequences
$$\begin{array}{ccccccccc}0& & ϵ𝒪_{Z_n}& & 𝒢& & 𝒪_{Z_n}& & 0\\ & & id& & D& & d& & & & \\ 0& & ϵ𝒪_{Z_n}& & & \stackrel{\varphi }{}& \mathrm{\Omega }_{Z_n/S_n}^1& & 0\end{array}$$
(1)
By definition $`𝒢=\{(f,g)𝒪_{Z_n};\varphi (f)=dg\}`$. We give a ring structure to $`𝒢`$ by $`(f,g)\times (f^{},g^{}):=(g^{}f+gf^{},gg^{})`$. By the inclusion $`ϵ𝒪_{Z_n}`$ we regard $`ϵ`$ as an element of $``$. Then $`𝒢`$ is an $`A_n[ϵ]/ϵ^2`$ algebra by defining $`ϵ.(f,g):=(gϵ,0)`$ for $`(f,g)𝒢`$. Now put $`𝒢=𝒪_{𝒵_n}`$. Then $`(Z_n,𝒪_{𝒵_n})`$ is an infinitesimal deformation of $`(Z_n,𝒪_{Z_n})`$ over $`S_n[ϵ]`$.
Remark. Another approach to prove (iii-a) is to use the cotangent complex and the associated tangent cohomology (cf. \[L-S\], \[Pa\], \[Bi, §. 5\]). Let $`L_{Z_n/S_n}^{}`$ be the cotangent complex for $`Z_nS_n`$. We put
$$T_{Z_n/S_n}^i:=\underset{¯}{\mathrm{Ext}}^i(L_{Z_n/S_n}^{},𝒪_{Z_n})$$
and
$$𝐓_{Z_n/S_n}^i:=\mathrm{𝐄𝐱𝐭}^i(L_{Z_n/S_n}^{},𝒪_{Z_n}).$$
By the definition of the cotangent complex, one has $`T_{Z_n/S_n}^0T_{Z_n/S_n}`$. Moreover, if $`Z`$ is reduced, then the argument of \[Ar 2, Prop. 6.1\] shows that $`T_{Z_n/S_n}^1\underset{¯}{\mathrm{Ext}}^1(\mathrm{\Omega }_{Z_n/S_n}^1,𝒪_{Z_n})`$. Since there is a natural map of complexes $`L_{Z_n/S_n}^{}\mathrm{\Omega }_{Z_n/S_n}^1`$, we get the maps
$$\underset{¯}{\mathrm{Ext}}^i:=\underset{¯}{\mathrm{Ext}}^i(\mathrm{\Omega }_{Z_n/S_n}^1,𝒪_{Z_n})T_{Z_n/S_n}^i$$
and
$$\mathrm{Ext}^i:=\mathrm{Ext}^i(\mathrm{\Omega }_{Z_n/S_n}^1,𝒪_{Z_n})𝐓_{Z_n/S_n}^i.$$
Therefore, we have a commutative diagram of local-to-global exact sequences:
$$\begin{array}{ccccccc}H^1(T_{Z_n/S_n})& & \mathrm{Ext}^1& & H^0(\underset{¯}{\mathrm{Ext}}^1)& & H^2(T_{Z_n/S_n})\\ & & & & & & & & \\ H^1(T_{Z_n/S_n}^0)& & 𝐓_{Z_n/S_n}^1& & H^0(T_{Z_n/S_n}^1)& & H^2(T_{Z_n/S_n}^0)\end{array}$$
(2)
By this diagram we see that $`\mathrm{Ext}^1(\mathrm{\Omega }_{Z_n/S_n}^1,𝒪_{Z_n})𝐓_{Z_n/S_n}^1`$ when $`Z`$ is reduced.
Next let us consider a locally trivial deformation $`Z_nS_n`$ of $`Z`$. Locally trivial deformations are controlled by $`H^1(Z,T_{Z_n/S_n})`$ instead of $`\mathrm{Ext}^1`$:
(iii-b): There is a one-to-one correspondence between $`H^1(Z,T_{Z_n/S_n})`$ and the set of equivalence classes of locally trivial deformations of $`Z_n`$ over $`S_n[ϵ]`$.
First let us take a Stein open cover $`\{V_i\}_{iI}`$ of $`Z_n`$. Then each locally trivial deformation $`𝒵_n`$ of $`Z_n`$ is constructed by patching $`V_i\times _{S_n}S_n[ϵ]`$ together in such a way that, over $`S_n`$, it gives the original $`Z_n`$. Now the situation is the same as the usual deformation of a non-singular variety. By the same argument as \[SGA 1\], Expose.III, 6, we have the above correspondence.
(iv) Assume that $`Z`$ is a Q-factorial projective symplectic variety with terminal singularities. By (ii), $`\mathrm{codim}(\mathrm{\Sigma }Z)4`$, where $`\mathrm{\Sigma }:=\mathrm{Sing}(Z)`$. We define $`A_n`$ and $`S_n`$ as in (iii). Suppose that $`Z_nS_n`$ is an infinitesimal deformation of $`Z`$ over $`S_n`$. Let $`U`$ be the regular part of $`Z`$. We put $`U_n:=Z_n|_U`$.
Lemma 1. The natural map
$$H^1(Z,T_{Z_n/S_n})\mathrm{Ext}^1(\mathrm{\Omega }_{Z_n/S_n}^1,𝒪_{Z_n})$$
is a bijection.
Proof. It suffices to show that the map
$$\mathrm{Ext}^1(\mathrm{\Omega }_{Z_n/S_n}^1,𝒪_{Z_n})H^0(Z,\underset{¯}{\mathrm{Ext}}^1(\mathrm{\Omega }_{Z_n/S_n}^1,𝒪_{Z_n}))$$
is the zero map. Take a point $`p\mathrm{\Sigma }`$. Let $`Z(p)`$ be a small Stein open neighborhood of $`pZ`$ and set $`U(p):=Z(p)U`$. We put $`Z_n(p):=Z_n|Z(p)`$ and $`U_n(p):=U_n|U(p)`$. We prove that
$$\mathrm{Ext}^1(\mathrm{\Omega }_{Z_n/S_n}^1,𝒪_{Z_n})H^0(Z_n(p),\underset{¯}{\mathrm{Ext}}^1(\mathrm{\Omega }_{Z_n(p)/S_n}^1,𝒪_{Z_n(p)}))$$
is the zero map. Since $`Z`$ is Cohen-Macaulay and $`\mathrm{codim}(\mathrm{\Sigma }Z)3`$, by \[Ko-Mo, 12.5.6\], one has
$$\mathrm{Ext}^1(\mathrm{\Omega }_{Z_n/S_n}^1,𝒪_{Z_n}))H^1(U_n,T_{U_n/S_n})$$
and
$$H^0(Z_n(p),\underset{¯}{\mathrm{Ext}}^1(\mathrm{\Omega }_{Z_n(p)/S_n}^1,𝒪_{Z_n(p)}))H^1(U_n(p),T_{U_n(p)/S_n}).$$
Since the symplectic 2-form $`\omega `$ extends to $`\omega _nH^0(U_n,\mathrm{\Omega }_{U_n/S_n}^2)`$ by \[Na\], we can identify $`T_{U_n/S_n}`$ with $`\mathrm{\Omega }_{U_n/S_n}^1`$ by this relative 2-form. Now we only have to prove that
$$H^1(U_n,\mathrm{\Omega }_{U_n/S_n}^1)H^1(U_n(p),\mathrm{\Omega }_{U_n(p)/S_n}^1)$$
is the zero map. Let us consider the Hodge spectral sequences
$$H^j(U_n,\mathrm{\Omega }_{U_n/S_n}^i)H^{i+j}(U,A_n)$$
and
$$H^j(U_n(p),\mathrm{\Omega }_{U_n(p)/S_n}^i)H^{i+j}(U(p),A_n),$$
where $`A_n`$ is the constant sheaf with coefficient $`A_n`$. In the latter spectral sequence, we have a decreasing filtration $`0F_{loc}^2F_{loc}^1F_{loc}^0=H^2(U(p),A_n)`$. By the depth argument, $`E_1^{0,1}:=H^1(U_n(p),𝒪_{U_n(p)})=0`$ because $`Z`$ is Cohen-Macaulay and $`\mathrm{codim}(\mathrm{\Sigma }Z)3`$. So there is an injection $`Gr_{F_{loc}}^1(H^2(U(p),A_n)H^1(U_n(p),\mathrm{\Omega }_{U_n(p)/S_n}^1)`$. In particular, we have a map $`F_{loc}^1H^1(U_n(p),\mathrm{\Omega }_{U_n/S_n}^1)`$. On the other hand, the first spectral sequence degenerates at $`E_1`$ terms for $`i+j=2`$ (cf. \[Na 1\], Lemma 2.7) because $`\mathrm{codim}(\mathrm{\Sigma }Z)4`$. Hence, for the decreasing filtration $`F^{}`$ of $`H^2(U,A_n)`$, $`Gr_F^1(H^2(U,A_n))=H^1(U_n,\mathrm{\Omega }_{U_n/S_n}^1)`$. In particular, we have a surjection $`F^1H^1(U_n,\mathrm{\Omega }_{U_n/S_n}^1)`$. The inclusion map $`U(p)U`$ induces a map of the second cohomologies, say $`\iota (p)`$
$$\iota (p):H^2(U,A_n)H^2(U(p),A_n).$$
This map preserves the filtrations $`\{F^{}\}`$ and $`\{F_{loc}^{}\}`$; hence it induces a map $`F^1F_{loc}^1`$. By the commutative diagram
$$F^1H^1(U_n,\mathrm{\Omega }_{U_n/S_n}^1)$$
$$$$
$$F_{loc}^1H^1(U_n(p),\mathrm{\Omega }_{U_n/S_n}^1)$$
we only have to prove that the map $`\iota (p)`$ is the zero map. Take a resolution $`\nu :\stackrel{~}{Z}Z`$ such that $`\nu ^1(U)U`$. Put $`\stackrel{~}{Z}(p):=\nu ^1(Z(p))`$. Let us consider the commutative diagram
$$H^2(\stackrel{~}{Z},A_n)H^2(U,A_n)$$
$$$$
$$H^2(\stackrel{~}{Z}(p),A_n)H^2(U(p),A_n).$$
Here note that $`H^2(,A_n)=H^2(,𝐂)_𝐂A_n`$. Since $`Z`$ is Q-factorial, $`\mathrm{im}[H^2(\stackrel{~}{Z},A_n)H^2(\stackrel{~}{Z}(p),A_n)]`$ is generated by $`[E_i]`$’s as an $`A_n`$ module by (i), where $`E_i`$ are exceptional divisors of $`\nu `$. Therefore, the composite
$$H^2(\stackrel{~}{Z},A_n)H^2(\stackrel{~}{Z}(p),A_n)H^2(U(p),A_n)$$
is the zero map. On the other hand, the map $`H^2(\stackrel{~}{Z},A_n)H^2(U,A_n)`$ is a surjection by an argument in \[Na, Proposition 9\](cf. (b), p.143) because $`\mathrm{codim}(\mathrm{\Sigma }Z)4`$. Now $`\iota (p)`$ is the zero map by the commutative diagram. Q.E.D.
Let us start the proof of Main Theorem. Let $`A`$ be an Artinian local $`𝐂`$ algebra and put $`S:=\mathrm{Spec}(A)`$. Recall that a deformation $`𝒵S`$ of $`Z`$ is locally trivial if, for any Stein open set $`VZ`$, $`𝒵|_VV\times S`$. We can introduce a locally trivial deformation functor of $`Z`$
$$D_{lt}:(\mathrm{Art})_𝐂(\mathrm{Set})$$
as a sub-functor of the usual deformation functor $`D`$ of $`Z`$. One can check that $`D_{lt}`$ has a pro-representable hull $`(R_{lt},𝒵_{R_{lt}})`$ in the sense of Schlessinger \[Sch\]. Here $`R_{lt}`$ is a complete local $`𝐂`$ algebra and $`𝒵_{R_{lt}}`$ is a certain projective system of infinitesimal deformations of $`Z`$. Moreover, by the definition of $`D_{lt}`$, one can construct $`R_{lt}`$ as a suitable quotient of the hull $`R`$ of the usual deformation functor $`D`$. $`D`$ has $`T^1`$-lifting property as noted in \[Na\] (for the relationship between $`T^1`$-lifting property and unobstructedness of deformations, see \[Ka\], \[Na 3\]). Hence, by Lemma 1 and (iii-b), we see that $`D_{lt}`$ also has $`T^1`$-lifting property. This implies that $`R`$ and $`R_{lt}`$ are both complete regular local ring over $`𝐂`$. Since their cotangent spaces coincide again by Lemma 1, we conclude that $`RR_{lt}`$. This implies that any flat deformation of $`Z`$ is locally trivial.
## 3 Applications
In this section we shall generalize Theorem 2.2 of \[Na 1\] to the case of Q-factorial terminalizations (cf. Theorem 1) and apply Main Theorem and Theorem 1 to get Corollaries 1 and 2. Propositions 1 and 2 are preliminaries for the proof of Theorem 1.
Let us recall some generalities on deformation theory. Let $`Y`$ be a compact complex space. Then it has been proved by \[Gr\], \[Do\], \[Pa\] and others that, there is a Kuranishi space $`0\mathrm{Def}(Y)`$ and a semi-universal deformation $`f:𝒴\mathrm{Def}(Y)`$ of $`Y`$ with $`f^1(0)=Y`$. Assume that $`h:YZ`$ is a (proper) map of compact complex spaces such that $`R^1h_{}𝒪_Y=0`$ and $`h_{}𝒪_Y=𝒪_Z`$. Let $`𝒴\mathrm{Def}(Y)`$ and $`𝒵\mathrm{Def}(Z)`$ be the semi-universal deformations. Then one can construct two maps $`h_{}:\mathrm{Def}(Y)\mathrm{Def}(Z)`$ and $`H_{}:𝒴𝒵`$ in such a way that the following diagram commutes (cf. \[Ko-Mo, 11.4\])
$$\begin{array}{ccc}𝒴& \stackrel{H_{}}{}& 𝒵\\ & & & & \\ \mathrm{Def}(Y)& \stackrel{h_{}}{}& \mathrm{Def}(Z)\end{array}$$
(3)
For the convenience of readers, we shall give a sketch of the construction of $`H_{}`$ and $`h_{}`$ according to \[Ko-Mo, 11.4\]. First note that $`H_{}`$ and $`h_{}`$ exist in the formal category by \[Wa\]. We want to construct them in the analytic category. Let $`GY\times Z`$ be the graph of $`h`$. We shall find a graph $`𝒢𝒴\times 𝒵`$ in such a way that $`𝒢(Y\times Z)=G`$. Let us consider the component $`D`$ of the relative Douady space of $`𝒴\times 𝒵/\mathrm{Def}(Y)\times \mathrm{Def}(Z)`$ containing $`[G]`$. The projection morphism $`D\mathrm{Def}(Y)`$ has a formal section defined by the formal contraction morphism around $`0\mathrm{Def}(Y)`$. Then, by \[Ar 1, 1.5\], one can find an analytic section. Note that, when $`h`$ is a birational map, $`(H_{})_t:Y_tZ_{h_{}(t)}`$ is a birational map for $`t\mathrm{Def}(Y)`$.
Proposition 1. Let $`Z`$ be a projective symplectic variety. Assume that $`h:YZ`$ is a Q-factorial terminalization, that is, $`h`$ is a birational morphism and $`Y`$ is a Q-factorial projective symplectic variety with terminal singularities. Then there is a deformation of $`h`$: $`𝒴_\mathrm{\Delta }\stackrel{\stackrel{~}{h}}{}𝒵_\mathrm{\Delta }\mathrm{\Delta }`$ over a disc such that for $`t\mathrm{\Delta }0`$, $`\stackrel{~}{h}_t:Y_tZ_t`$ is an isomorphism.
Proof. Denote by $`(T,0)`$ the Kuranishi space of $`Y`$ and let $`p:𝒴T`$ be the semi-universal family over $`T`$. By \[Na 1\], Proposition (2.5), $`T`$ is smooth. Take a projective resolution $`\pi :\stackrel{~}{𝒴}𝒴`$ in such a way that $`\pi `$ is an isomorphism outside $`\mathrm{Sing}(𝒴)`$. For a general $`tT`$, $`\pi _t:\stackrel{~}{Y}_tY_t`$ is a resolution. Moreover, $`\pi _t^1(U_t)U_t`$ for the regular locus $`U_t`$ of $`Y_t`$, since $`U_t\mathrm{Sing}(𝒴)=\mathrm{}`$. Note that $`\stackrel{~}{Y}_t`$ is not, in general, a projective variety, but a Kähler manifold. We shall prove that one can deform $`Y_t`$ further to a variety $`Y_t^{}`$ so that it contains no curves. Since there is a map of Kuranishi spaces $`h_{}:T\mathrm{Def}(Z)`$ as noted above, one has a birational morphism $`Y_t^{}Z_{h_{}(t^{})}`$. But, $`Y_t^{}`$ does not have any curve; hence this birational morphism should be an isomorphism (cf. \[Na 1\], claim 3, p. 618). Since $`Y_t`$ is a small deformation of a symplectic variety $`Y`$, $`Y_t`$ is again symplectic by \[Na\]. A symplectic singularity is a canonical singularity (cf. \[Be\]). A canonical singularity is a rational singularity by \[El\]. Therefore, $`Y_t`$ has only rational singularities. Since $`Y_t`$ has only rational singularities, the natural map $`H^2(Y_t,𝐂)H^2(\stackrel{~}{Y}_t,𝐂)`$ is an injection of mixed Hodge structures (\[De\]). Hence, by the strict compatibility of weight filtrations, $`H^2(Y_t,𝐂)`$ admits a pure Hodge structure of weight 2. Let
$$H^2(Y_t,𝐂)=H^{2,0}(Y_t)H^{1,1}(Y_t)H^{0,2}(Y_t)$$
be the Hodge decomposition. Recall again that there is an injection of pure Hodge structures $`H^2(Y_t,𝐂)H^2(\stackrel{~}{Y}_t,𝐂)`$. Put $`n:=dimY`$ and consider the perfect pairing
$$<,>:H^{2n2}(\stackrel{~}{Y}_t,𝐂)\times H^2(\stackrel{~}{Y}_t,𝐂)H^{2n}(\stackrel{~}{Y}_t,𝐂)=𝐂.$$
Set
$$V_t:=H^2(Y_t,𝐂)^{}H^{2n2}(\stackrel{~}{Y}_t,𝐂).$$
Here we regard $`H^2(Y_t,𝐂)`$ as a subspace of $`H^2(\stackrel{~}{Y}_t,𝐂)`$.
Claim 1. Every element of $`V_t`$ is of type $`(n1,n1)`$.
Proof. If we put
$$V_t^{n1,n1}:=V_tH^{n1}(\stackrel{~}{Y}_t,\mathrm{\Omega }_{\stackrel{~}{Y}_t}^{n1}),$$
then we see that $`V_t^{n1,n1}=H^{1,1}(Y_t)^{}H^{n1}(\stackrel{~}{Y}_t,\mathrm{\Omega }_{\stackrel{~}{Y}_t}^{n1})`$ for the Serre pairing
$$<,>:H^{n1}(\stackrel{~}{Y}_t,\mathrm{\Omega }_{\stackrel{~}{Y}_t}^{n1})\times H^1(\stackrel{~}{Y}_t,\mathrm{\Omega }_{\stackrel{~}{Y}_t}^1)H^n(\stackrel{~}{Y}_t,\mathrm{\Omega }_{\stackrel{~}{Y}_t}^n)=𝐂.$$
Since $`h^{n1,n1}(\stackrel{~}{Y}_t)=h^{1,1}(\stackrel{~}{Y}_t)`$, $`dimV_t^{n1,n1}=h^{1,1}(\stackrel{~}{Y}_t)h^{1,1}(Y_t)`$. On the other hand, $`dimV_t=b_2(\stackrel{~}{Y}_t)b_2(Y_t)`$. We shall prove that $`dimV_t=dimV_t^{n1,n1}`$. In order to do this, it suffices to show that all elements of $`H^2(\stackrel{~}{Y}_t)`$ of type (2,0) and of type (0,2) are mapped to zero by the map $`\varphi :H^2(\stackrel{~}{Y_t},𝐂)H^0(Y_t,R^2\pi _{t}^{}{}_{}{}^{}𝐂)`$. In fact, since the map $`H^2(Y_t)H^2(\stackrel{~}{Y}_t)`$ is an injection of pure Hodge structures of weight 2, if these are proved, then the injection is actually an isomorphism at (2,0) and (0,2) parts. By the conjugation, it suffices to show that $`\varphi (\alpha )=0`$ for an element $`\alpha H^2(\stackrel{~}{Y}_t)`$ of type (2,0). We shall prove that, for every point $`xY_t`$, $`\varphi (\alpha )_x=0`$ in $`(R^2(\pi _t)_{}𝐂)_x`$. Let $`\nu :W\stackrel{~}{Y}_t`$ be a projective birational morphism such that $`W`$ is smooth and $`D:=(\pi _t\nu )^1(x)`$ is a simple normal crossing divisor of $`W`$. Put $`h:=\pi _t\nu `$. Since $`R^1\nu _{}𝐂=0`$, $`R^2(\pi _t)_{}𝐂`$ injects to $`R^2h_{}𝐂`$; hence we have to check that $`\alpha `$ is sent to zero by the composite
$$H^2(\stackrel{~}{Y}_t,𝐂)H^2(W,𝐂)(R^2h_{}𝐂)_x(=H^2(D,𝐂)).$$
We call this composite $`\psi `$. Then $`\psi `$ preserves Hodge filtrations of $`H^2(\stackrel{~}{Y}_t)`$ and $`H^2(D)`$. Hence it induces $`Gr_F^2(H^2(\stackrel{~}{Y}_t))Gr_F^2(H^2(D))H^2(D)`$. Since $`\alpha Gr_F^2(H^2(\stackrel{~}{Y}_t))`$, $`\psi (\alpha )Gr_F^2(H^2(D))`$. But, since $`Y_t`$ has rational singularities at $`x`$, $`Gr_F^2(H^2(D))=H^0(D,\widehat{\mathrm{\Omega }}_D^2)=0`$ by \[Na 1, Lemma (1.2)\], where $`\widehat{\mathrm{\Omega }}_D^2:=\mathrm{\Omega }_D^2/(\mathrm{torsion})`$. Therefore, $`dimV_t=dimV_t^{n1,n1}`$. Q.E.D.
The families
$$\stackrel{~}{𝒴}𝒴T$$
induces the Kodaira-Spencer maps
$$\rho _{\stackrel{~}{Y}_t}:T_t(T)H^1(\stackrel{~}{Y}_t,\mathrm{\Theta }_{\stackrel{~}{Y}_t}),$$
$$\rho _{Y_t}:T_t(T)\mathrm{Ext}^1(\mathrm{\Omega }_{Y_t}^1,𝒪_{Y_t}).$$
Since $`𝒴T`$ is versal at $`tT`$, $`\rho _{Y_t}`$ is surjective. Define a map
$$d:H^1(\stackrel{~}{Y}_t,\mathrm{\Theta }_{\stackrel{~}{Y}_t})\mathrm{Ext}^1(\mathrm{\Omega }_{Y_t}^1,𝒪_{Y_t})$$
as the differential of the map $`\mathrm{Def}(\stackrel{~}{Y}_t)T`$. We then have $`d\rho _{\stackrel{~}{Y}_t}=\rho _{Y_t}`$. Let $`U_tY_t`$ be the regular locus of $`Y_t`$. Then $`\mathrm{Ext}^1(\mathrm{\Omega }_{Y_t}^1,𝒪_{Y_t})H^1(U_t,\mathrm{\Theta }_{U_t})`$. By this identification, the map $`d`$ coincides with the restriction map $`H^1(\stackrel{~}{Y}_t,\mathrm{\Theta }_{\stackrel{~}{Y}_t})H^1(U_t,\mathrm{\Theta }_{U_t})`$. A symplectic 2-form on $`Y_t`$ extends to a holomorphic 2-form on $`\stackrel{~}{Y}_t`$. By this 2-form, we have a map $`H^1(\mathrm{\Theta }_{\stackrel{~}{Y}_t})H^1(\mathrm{\Omega }_{\stackrel{~}{Y}_t}^1)`$. We put
$$T_1:=\mathrm{im}[T_t(T)H^1(\stackrel{~}{Y}_t,\mathrm{\Theta }_{\stackrel{~}{Y}_t})],$$
$$T_2:=\mathrm{im}[T_t(T)H^1(\stackrel{~}{Y}_t,\mathrm{\Omega }_{\stackrel{~}{Y}_t}^1)].$$
Claim 2.
$$dimT_2dimH^1(U_t,\mathrm{\Omega }_{U_t}^1).$$
Proof. Let us consider the composite
$$T_t(T)H^1(\stackrel{~}{Y_t},\mathrm{\Omega }_{\stackrel{~}{Y}_t}^1)H^1(U_t,\mathrm{\Omega }_{U_t}^1).$$
By the symplectic 2-form, we identify $`H^1(U_t,\mathrm{\Theta }_{U_t})`$ with $`H^1(U_t,\mathrm{\Omega }_{U_t}^1)`$. Then the composite above coincides with the Kodaira-Spencer map $`\rho _{Y_t}`$, which is surjective.
Look at the pairing
$$(,):H^1(\stackrel{~}{Y}_t,\mathrm{\Theta }_{\stackrel{~}{Y}_t})\times H^{n1}(\stackrel{~}{Y}_t,\mathrm{\Omega }_{\stackrel{~}{Y}_t}^{n1})H^n(\stackrel{~}{Y}_t,\mathrm{\Omega }_{\stackrel{~}{Y}_t}^{n2})$$
Claim 3. In the pairing above, $`(T_1,V_t)=0`$.
Proof. For $`\zeta H^1(\stackrel{~}{Y}_t,\mathrm{\Theta }_{\stackrel{~}{Y}_t})`$, take the corresponding infinitesimal deformation $`\stackrel{~}{Y}_{t,ϵ}\mathrm{Spec}𝐂[ϵ]`$. Choose an element $`lH^{2n2}(\stackrel{~}{Y}_t,𝐂)`$ of type $`(n1,n1)`$. Then $`(\zeta ,l)`$ is an obstruction for the class $`l`$ to remain of type $`(n1,n1)`$ under the infinitesimal deformation. Now the claim follows from Claim 1 and the fact that $`R^2p_{}𝐂_𝒴`$ is a local system on $`T`$ (cf. \[Na\]<sup>1</sup><sup>1</sup>1In \[Na\], this has been proved under ceratin additional conditions: $`h^0(Y_{reg},\mathrm{\Omega }_{Y_{reg}}^2)=1`$, $`h^1(Y,𝒪_Y)=0`$. However, these conditions are not essential. Here we give a more direct proof to this fact by using our Main Theorem. Since $`𝒴T`$ is proper, if we choose $`T`$ small enough, then there is an open covering $`\{𝒱_i\}_{iI}`$ of $`𝒴`$ with (analytic) trivializations over $`T`$: $`𝒱_iV_i\times T`$. Here $`V_i:=𝒴𝒱_i`$. Let us take a canonical stratification $`Y=_jY_j`$ of $`Y`$ into locally closed smooth subsets such that $`\overline{Y}_j=\mathrm{Sing}(\overline{Y}_{j1})`$ and $`Y_0=Y_{reg}`$. Fix a nowhere vanishing $`C^{\mathrm{}}`$-vector field $`\zeta `$ of $`T`$. One can lift $`\zeta `$ to a $`C^{\mathrm{}}`$-vector field $`\zeta _i`$ of $`V_i\times T`$ in such a way that $`\zeta _i`$ induces a vector field of each $`(Y_jV_i)\times T`$ for each stratum $`Y_j`$. By using the partition of unity, one can glue together $`\{\zeta _i\}`$ and get a globally defined vector field $`\stackrel{~}{\zeta }`$ of $`𝒴`$. This vector field gives us a $`C^{\mathrm{}}`$-trivialization of $`𝒴`$ over $`T`$. Hence $`R^ip_{}𝐂`$ are constant sheaves.).
Since $`T_2`$ is the image of $`T_1`$ by the map $`H^1(\mathrm{\Theta }_{\stackrel{~}{Y}_t})H^1(\mathrm{\Omega }_{\stackrel{~}{Y}_t}^1)`$, we have $`<T_2,V_t>=0`$ by the commutative diagram of the pairing maps:
$$(,):H^1(\stackrel{~}{Y}_t,\mathrm{\Theta }_{\stackrel{~}{Y}_t})\times H^{n1}(\stackrel{~}{Y}_t,\mathrm{\Omega }_{\stackrel{~}{Y}_t}^{n1})H^n(\stackrel{~}{Y}_t,\mathrm{\Omega }_{\stackrel{~}{Y}_t}^{n2})$$
$$$$
$$<,>:H^1(\stackrel{~}{Y}_t,\mathrm{\Omega }_{\stackrel{~}{Y}_t}^1)\times H^{n1}(\stackrel{~}{Y}_t,\mathrm{\Omega }_{\stackrel{~}{Y}_t}^{n1})H^n(\stackrel{~}{Y}_t,\mathrm{\Omega }_{\stackrel{~}{Y}_t}^n).$$
Note that $`\mathrm{codim}(V_tH^{n1,n1}(\stackrel{~}{Y}_t))=h^{1,1}(Y_t)`$. On the other hand, since $`h^{1,1}(Y_t)=h^1(U_t,\mathrm{\Omega }_{U_t}^1)`$(cf. \[Na\]<sup>2</sup><sup>2</sup>2By the argument of the footenote (1) above, we see that $`h^2(Y)=h^2(Y_t)`$ and $`h^2(U)=h^2(U_t)`$ with $`U:=Y_{reg}`$. Since $`Y`$ is Q-factorial, one has $`H^2(Y)H^2(U)`$; hence $`H^2(Y_t)H^2(U_t)`$. Now we get the result by the footnote (1) in p. 21 of the eprint version of \[Na\] (math.AG/0010114).), we see, by Claim 2, that $`T_2=(V_t)^{}`$ with respect to the pairing $`<,>`$ above.
Now let $`C`$ be a connected curve on $`\stackrel{~}{Y}_t`$ such that $`\pi _t(C)`$ is not a point. Since $`[C]V_t`$, we see that $`[C]^{}T_2`$ is a codimension 1 subspace of $`T_2`$ by the perfectness of $`<,>`$. In particular, we can find an element $`\zeta T_1`$ such that $`(\zeta ,[C])0`$. The $`(n1,n1)`$ classes of curves on $`\stackrel{~}{Y}_t`$ constitute a countable subset of $`H^{n1,n1}(\stackrel{~}{Y}_t)`$. Hence, if we take a generic deformation $`\stackrel{~}{Y}_t^{}`$ of $`\stackrel{~}{Y}_t`$, then any effective 1-cycle on $`\stackrel{~}{Y}_t`$ which is not contaied in a fiber of $`\pi _t`$ cannot deform sideways. Since $`\stackrel{~}{𝒴}T`$ is a Kaehler morphism, every irreducible component of the Douady space $`D(\stackrel{~}{𝒴}/T)`$ parametrizing curves is proper over $`T`$ (\[Fu 2\]). Therefore, all components dominating $`T`$ parametrize the curves which are contracted to points by $`\stackrel{~}{𝒴}𝒴`$. By the countability of the Douady space (\[Fu 3\]), we now see that, for a generic deformation $`\stackrel{~}{Y}_t^{}`$ of $`\stackrel{~}{Y}_t`$, any holomorphic curve $`C`$ on $`\stackrel{~}{Y}_t^{}`$ is $`\pi _t^{}`$-exceptional for $`\pi _t^{}:\stackrel{~}{Y}_t^{}Y_t^{}`$, that is, $`\pi _t^{}(C)`$ is a point. Then, the variety $`Y_t^{}`$ has no holomorphic curves. In fact, if there is a curve $`C`$, then one can find a curve $`D`$ on $`\stackrel{~}{Y}_t^{}`$ such that $`\pi _t^{}(D)=C`$ by using the Chow lemma \[Hi\], which is a contradiction.
Let $`h:YZ`$ be a Q-factorial terminalization of a projective symplectic variety $`Z`$. We put $`\mathrm{\Sigma }:=\mathrm{Sing}(Z)`$. There is a closed subset $`\mathrm{\Sigma }_0`$ of $`\mathrm{\Sigma }`$ with $`\mathrm{codim}(\mathrm{\Sigma }_0Z)4`$ such that $`Z`$ is locally a trivial family of a rational double point along $`\mathrm{\Sigma }\mathrm{\Sigma }_0`$. In the remainder, we put $`U:=Z\mathrm{\Sigma }_0`$ and $`V:=Y\mathrm{Sing}(Y)`$. Note that $`h^1(U)V`$. We denote by $`n`$ the dimension of $`Z`$. In this situation, we have the following generalization of \[Na 1\], Proposition (2.1).
Proposition 2<sup>3</sup><sup>3</sup>3Q-factoriality of $`Y`$ is not necessary for this proposition. There is a commutative diagram
$$\mathrm{Ext}^1(\mathrm{\Omega }_Y^1,𝒪_Y)H^1(h^1(U),\mathrm{\Theta }_{h^1(U)})$$
$$$$
$$\mathrm{Ext}^1(\mathrm{\Omega }_Z^1,𝒪_Z)\mathrm{Ext}^1(\mathrm{\Omega }_U^1,𝒪_U),$$
where horizontal maps are both isomorphisms.
Proof. In the diagram, each space corresponds to the set of first order deformations of $`Y`$, $`h^1(U)`$, $`Z`$ or $`U`$, respectively. A first order deformation of $`Y`$ (resp. $`h^1(U)`$) induces that of $`Z`$ (resp. $`U`$)(cf. \[Na 1\], p.614). Each vertical map is nothing but this correspondence. On the other hand, the horizontal maps are natural ones induced by the restriction. The second horizontal map (at the bottom) is an isomorphism by \[Ko-Mo, 12.5.6\] because $`\mathrm{codim}(\mathrm{\Sigma }_0Z)3`$ and $`Z`$ is Cohen-Macaulay. By the same reason, $`\mathrm{Ext}^1(\mathrm{\Omega }_Y^1,𝒪_Y)H^1(V,\mathrm{\Theta }_V).`$ We only have to prove that the first horizontal map is an isomorphism. Set $`F:=Yh^1(U)`$ and $`F^0:=Vh^1(U)`$. Note that $`dimFn2`$ since $`h`$ is a crepant partial resolution of a symplectic variety (\[Na 1\], Corollary (1.15); see also \[Na 2\], footnote, p.1). Let us consider the exact sequence of cohomology with coefficient $`𝐂`$
$$H_{F^0}^2(V)H^2(V)H^2(h^1(U))H_{F^0}^3(V).$$
Let $`j:h^1(U)V`$ and $`i:F^0V`$ be the inclusion maps. Then the dual of this sequence coincides with the exact sequence of cohomology with compact support:
$$H_c^{2n3}(F^0)H_c^{2n2}(h^1(U))H_c^{2n2}(V)H_c^{2n2}(F^0)$$
induced by the exact sequence (cf. \[Ha,II, Ex.1.19(c)\])
$$0j_!(𝐂_{h^1(U)})𝐂_Vi_{}(𝐂_{F^0})0.$$
Here we have $`H_c^{2n3}(F^0)=H_c^{2n2}(F^0)=0`$. In fact, the first terms and the third terms of the next exact sequences vanish (cf.\[L\]) since $`dim_𝐂Fn2`$ and $`dim_𝐂(FF^0)n3`$:
$$H^{2n3}(FF^0)H_c^{2n2}(F^0)H^{2n2}(F)$$
and
$$H^{2n4}(FF^0)H_c^{2n3}(F^0)H^{2n3}(F).$$
As a consequence, the map
$$H^2(V)H^2(h^1(U))$$
is an isomorphism. Moreover, this is a morphism of mixed Hodge structures. Since $`Y`$ is a symplectic variety with terminal singularities, $`\mathrm{codim}(\mathrm{Sing}(Y)Y)4`$; hence $`H^2(V)`$ is equipped with a pure Hodge structure of weight 2 with $`\mathrm{Gr}_F^1(H^2(V))=H^1(\mathrm{\Omega }_V^1)`$ (cf. e-print version math.AG/0010114 of \[Na\], footnote at p.21). Let us examine the mixed Hodge structure on $`H^2(h^1(U))`$. Let $`\nu :\stackrel{~}{Y}Y`$ be a projective birational morphism from a non-singular variety $`\stackrel{~}{Y}`$ to $`Y`$ such that $`\nu ^1(h^1(U))h^1(U)`$ and such that $`\mathrm{Exc}(h\nu )`$ is a normal crossing variety of $`\stackrel{~}{Y}`$. We put $`E^{}:=\stackrel{~}{Y}(h\nu )^1(U)`$. By \[De\], $`\mathrm{Gr}_F^1(H^2(h^1(U)))=H^1(\stackrel{~}{Y},\mathrm{\Omega }_{\stackrel{~}{Y}}^1(\mathrm{log}E^{}))`$. Let us consider the exact sequence of local cohomology
$$H_E^{}^1(\mathrm{\Omega }_{\stackrel{~}{Y}}^1(\mathrm{log}E^{}))H^1(\mathrm{\Omega }_{\stackrel{~}{Y}}^1(\mathrm{log}E^{}))H^1(h^1(U),\mathrm{\Omega }_{h^1(U)}^1)H_E^{}^2(\mathrm{\Omega }_{\stackrel{~}{Y}}^1(\mathrm{log}E^{})).$$
The first term and the fourth term both vanish. This follows from the same argument as the proof of \[Na 1\], Proposition (2.1): by \[Na 1, claim 2, p.616\], $`R^k(h\nu )_{}\mathrm{\Omega }_{\stackrel{~}{Y}}^{n1}(\mathrm{log}E^{})(E^{}))=0`$ for $`k2`$. Then by taking the duals of the first term and the fourth term, we get the conclusion since $`\mathrm{codim}(\mathrm{\Sigma }_0Z)4`$ (cf. \[Na 1, p.615\]). As a consequence, we see that
$$\mathrm{Gr}_F^1(H^2(h^1(U))=H^1(h^1(U),\mathrm{\Omega }_{h^1(U)}^1).$$
Since the natural map
$$H^1(V,\mathrm{\Omega }_V^1)H^1(h^1(U),\mathrm{\Omega }_{h^1(U)}^1)$$
coincides with the map
$$\mathrm{Gr}_F^1(H^2(V))\mathrm{Gr}_F^1(H^2(h^1(U))),$$
we conclude that it is an isomorphism. Finally, by the symplectic 2-form on $`V`$, this map is identified with
$$H^1(V,\mathrm{\Theta }_V)H^1(h^1(U),\mathrm{\Theta }_{h^1(U)}).$$
Theorem 1. Let $`h:YZ`$ be a Q-factorial terminalization of a projective symplectic variety $`Z`$. Then the Kuranishi spaces $`\mathrm{Def}(Y)`$ and $`\mathrm{Def}(Z)`$ are both smooth of the same dimension. The natural map (cf. \[Na 1\]) $`h_{}:\mathrm{Def}(Y)\mathrm{Def}(Z)`$ is a finite covering. Moreover, $`Z`$ has a flat deformation to a (non-projective) symplectic variety $`Z_t`$ which is, at the same time, a small deformation of $`Y`$.
Proof. When $`Y`$ is non-singular, it is nothing but Theorem (2.2) of \[Na 1\]. In a general case, we modify the proof of \[Na 1\] in the following way. First, Corollary (1.5) of \[Na 1\] holds true when $`\stackrel{~}{X}`$ is a Q-factorial terminalization of $`X`$ (here the notation being the same as \[Na 1\]). Next, Proposition (1.6) of \[Na 1\] follows from a more general fact recently proved in \[K\]. Finally, we have to replace Proposition (2.1) of \[Na 1\] by Proposition 2 and replace Claim 3 of \[Na 1\](p.618) by Proposition 1.
Corollary 1. Let $`h:YZ`$ and $`h^{}:Y^{}Z`$ be two Q-factorial terminalizations of a projective symplectic variety $`Z`$. Assume that $`Y`$ is non-singular. Then $`Y^{}`$ is also non-singular.
Proof. Since $`h`$ is a symplectic resolution, the Kuranishi space $`\mathrm{Def}(Z)`$ of $`Z`$ is non-singular by \[Na 1, Theorem (2.2)\]. Moreover, a general point of $`\mathrm{Def}(Z)`$ parametrizes a non-singular variety. Assume now that $`Y^{}`$ is singular. By Proposition 1, there is a point $`t\mathrm{Def}(Z)`$ parametrizing a suitable deformation $`Y_t^{}`$ of $`Y^{}`$. By Main Theorem, $`Y^{}`$ does not have any smoothing via deformation; hence $`Y_t^{}`$ does not have, too. This contradicts the fact that a general point of $`\mathrm{Def}(Z)`$ parametrizes a non-singular variety.
Corollary 2. Let $`Z`$ be a projective symplectic variety. Assume that $`Z`$ has a Q-factorial terminalization $`h:YZ`$. Then the following are equivalent.
(1) $`Z`$ is smoothable by a flat deformation;
(2) $`Y`$ is non-singular.
Proof. (1) $``$ (2): Suppose that $`Y`$ is singular. Take a small deformation $`Y_t`$ of $`Y`$. Then $`Y_t`$ is not smoothable by any flat deformation by Main Theorem. Applying Proposition 1 to $`h:YZ`$, we see that $`Z_{h_{}(t)}(Y_t)`$ is also non-smoothable. But, by Theorem 1, $`\mathrm{Def}(Z)`$ is non-singular and hence irreducible. This contradicts the assumption (1).
(2) $``$ (1): This is nothing but Proposition 1.
Yoshinori Namikawa
Departement of Mathematics, Graduate School of Science, Osaka University, JAPAN
namikawa@math.sci.osaka-u.ac.jp |
warning/0506/cond-mat0506759.html | ar5iv | text | # Two-component sandpile model : self-organized criticality of the second kind
## Abstract
Two-component sandpile models are investigated numerically and theoretically. Monte Calro simulations are performed to show that probability distribution functions of avalanche size and lifetime obey power laws whose exponents are approximately equal to $`1.5`$ and $`2.0`$ and the system exhibits SOC. A mean-field theory is developed to discuss the essence of the processes. We find that two-component models approach a steady critical state belonging to a different universality class from that of one-component models. Conservation of two kinds of sands at local toppling causes an infinite number of stable states which substitute for artificial boundary dissipation. Among two control parameters appearing in one-component models, therefore, a rate constant of dissipation is removed in two-component models. It is concluded that the more conserved quantities result in the less control parameters and a novel class of SOC.
Stochastic processes, Sandpile models, Self-organized criticality, (Power-law distribution)
According to a vast amount of studies on fractals, $`1/f`$ noise, and so on, it has been reported that a lot of nonequiliblium systems in nature exhibit power laws. However it has not been deeply understood why a huge variety of distributions spontaneously go to power laws. Self-organized criticality(SOC) proposed by Bak et al.BTW1987 ; BTW1988 ; Jensen ; Sornette has provided one reasonable understanding of the emergence of power laws through sandpile models. The models are cellular automata evolving in threshold dynamics. The domino effect induces a sequence of topplings called an avalanche, whose distributions of magnitude and lifetime obey power laws without tuning any parameters. A number of works on SOC have been done numerically. Several analytical results have also been obtained. DharDhar1990 ; Dhar1999 has found some exact solutions in abelian sandpile models and has shed light on the statics of the models, such as the number of total recurrent states and hight correlation functions. However, the dynamics of sandpile models is still obscure, that is, the power-law distributions of avalanches still have been unable to be derived analytically.
Though sandpile models seem to have no control parameters by appearances, Vespegnani and ZapperiVZ1997 ; VZ1998 have found that the models do have hidden control parameters: a slow dissipation rate and a slower addition rate of sands. To explain this, they have introduced a mean-field theory of rate equations,
$`{\displaystyle \frac{}{t}}\rho _\kappa =f_\kappa (\rho _a,\rho _c,\rho _s),\kappa =a,c,s`$ (1)
where $`\rho _a`$, $`\rho _c`$, and $`\rho _s`$ are the probability densities of active, critical, and stable states, respectively. Imposing the conservation law of the number of sands in a local toppling rule, it is concluded that the models become critical in the double limit
$`ϵ,\delta 0,\rho _a=\delta /ϵ0,`$ (2)
where $`ϵ`$ is a dissipation rate and $`\delta `$ an addition rate. In other words, $`\rho _a`$ is an order parameter and $`ϵ`$, $`\delta `$ are control parameters. SOC is achieved by rough tuning of the ”unapparent” parameters $`ϵ`$, $`\delta `$ around zero. Most of SOC modelsManna1991 ; DS1992 belong to this ”universality” class. We call this class SOC of the first kind. It seems interesting whether other kind of SOC exists, or equivalently, whether the condition of the double limit (2) are strictly the necessary condition of SOC? In this letter, we give an answer to these questions.
As mentioned above, the conservation law of the local toppling rule plays a key role for SOC. Tsuchiya and KatoriTK2000 have proved rigorously that SOC breaks down in abelian sandpile models with nonconservative toppling of sands. Here, we pay attention to the number of conservation laws. Does anything new happen if the number of conservation laws is increased? How about the robustness of SOC? To answer these questions, we consider two-component sandpile models which deal with two kinds of sands. This means the models have two kinds of conservation laws. Firstly, numerical simulations are carried out to check power-law behavior. Secondly, a mean-field theory is developed to examine the essence of the processes.
A two-component sandpile model is a cellular automaton defined on a regular lattice. A pair of two non-negative integers $`𝒉(𝒙)(i,j)`$ is assigned to each site $`𝒙`$ on the lattice, where $`i,j(0)`$ denote the numbers of sand $`A`$ and $`B`$. At each time step, one unit of sand $`A`$ or $`B`$ is added to a randomly chosen site at a relative ratio $`(0<)r_{AB}(<1)`$ of sands $`A`$ to sands $`B`$. Several types of toppling rule can be adopted. For example, a toppling occurs when (a) $`ii_{th}`$ or $`jj_{th}`$, (b) $`ii_{th}`$ and $`jj_{th}`$, or (c) $`i+jk_{th}`$ $`(i,j>0)`$ where $`i_{th}`$, $`j_{th}`$, and $`k_{th}`$ are certain threshold values. If the number of sands at one site reaches threshold, sands $`A`$ and $`B`$ on the site topple one by one to randomly chosen nearest neighbor sites until the number of sands at the site becomes less than threshold. In these rules, a sand $`A`$ and a sand $`B`$ topple jointly. For example, if a site has $`n`$ grains of sands $`A`$ and no sands $`B`$, and then $`m`$ ($`n`$) new grains of sands $`B`$ are added, it ends up with $`nm`$ grains of sands A and no sands $`B`$. These rules satisfy two kinds of conservation laws of local toppling. A series of topplings (an avalanche) continue unless the numbers of sands at all sites become less than threshold. Among three rules of toppling, the and rule (b) is the most interesting because of the presence of an infinite number of stable states. In the rule (b), stable states $`𝒉_s`$ and unstable states $`𝒉_u`$ are defined by
$`𝒉_s`$ $`=`$ $`\{(i,j)|\mathrm{\hspace{0.33em}0}i<i_{th}\text{or}\mathrm{\hspace{0.33em}0}j<j_{th}\},`$ (3)
$`𝒉_u`$ $`=`$ $`\{(i,j)|ii_{th}\text{and}jj_{th}\}.`$ (4)
and there are infinite stable and unstable states as illustrated in Fig. 1.
These stable states can absorb all the added sands. This means avalanches can stop without introducing artificial boundary dissipation, which is indispensable in one-component models. Rules (a) and (c) are not capable of removing boundary dissipation because the number of stable states is finite under these rules. Hereafter, we consider only the rule (b) and treat threshold $`i_{th}=j_{th}=1`$ and periodic boundary conditions.
Distribution functions of avalanche size $`S`$ and lifetime $`T`$ are calculated numerically, where $`S`$ is defined by the total number of topplings in one avalanche and $`T`$ by the total time steps of simultaneous updates of topplings. Cumulative distribution functions (CDFs) of avalanche size $`D(S)`$ and lifetime $`D(T)`$ in $`1`$dimensional lattices are shown in Figs. 2.
We find that the CDFs obey power laws $`D(S)S^\tau `$ and $`D(T)T^\alpha `$ where the exponents are approximately $`\tau 0.5`$ and $`\alpha 1.0`$. Equivalently, probability distribution functions (PDFs) are represented as $`P(S)S^{(1+\tau )}`$ and $`P(T)S^{(1+\alpha )}`$. These exponents $`1+\tau 1.5`$ and $`1+\alpha 2.0`$ are very close to the critical exponents of branching processesHarris . We carried out simulations in $`2`$ and $`3`$dimension and obtained almost the same CDFs as those in $`1`$dimension. These results imply the model has meanfield-like characteristicsAlstrom1988 . It is also found that varying $`r_{AB}`$ from $`0.5`$ to $`0.1`$ hardly affects power-law tails. This means the model ends up with going to SOC states robustly in the long time limit without tuning $`r_{AB}`$. It becomes evident that the two-component sandpile models exhibits SOC. It should be noted that because of the absence of boundary dissipation, the power-law tails become extended infinitely as the number of added sands $`N`$ increases even in a finite system.
Next, we construct a mean-field theory of two-component sandpile models and compare with that of one-component modelsVZ1997 ; VZ1998 . We introduce the sets of probability densities of stable states $`𝑿`$ and unstable states $`𝒁`$ as
$`𝑿`$ $`=`$ $`\{X_{(i,j)}|\mathrm{\hspace{0.33em}0}i<1\text{or}\mathrm{\hspace{0.33em}0}j<1\},`$ (5)
$`𝒁`$ $`=`$ $`\{Z_{(i,j)}|i1\text{and}j1\}.`$ (6)
Rate equations for $`𝑿`$ and $`𝒁`$ are given by an infinite system of first-order nonlinear ordinary differential equations. However, it is difficult to handle the infinite degrees of freedom. In order to avoid the difficulty, we define the following reduced six variables
$`𝑿^{\mathbf{}}=\{X_0,X_A,X_B\},𝒁^{\mathbf{}}=\{Z_0,Z_A,Z_B\}`$ (7)
with
$`X_0X_{(0,0)},Z_0{\displaystyle \underset{i=1}{\overset{\mathrm{}}{}}}Z_{(i,i)},`$
$`X_A{\displaystyle \underset{i=1}{\overset{\mathrm{}}{}}}X_{(i,0)},Z_A{\displaystyle \underset{i=2}{\overset{\mathrm{}}{}}}{\displaystyle \underset{j=1}{\overset{\mathrm{}}{}}}Z_{(i,j)},`$
$`X_B{\displaystyle \underset{j=1}{\overset{\mathrm{}}{}}}X_{(0,j)},Z_B{\displaystyle \underset{i=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{j=2}{\overset{\mathrm{}}{}}}Z_{(i,j)}.`$
For obtaining a closed set of equations for $`𝑿^{\mathbf{}}`$ and $`𝒁^{\mathbf{}}`$, we assume Poisson distributions for $`X_{(i,0)}`$ and $`X_{(0,j)}`$ as
$`{\displaystyle \frac{X_{(i,0)}}{X_A}}`$ $`=`$ $`{\displaystyle \frac{\mu _A^{i1}}{(i1)!}}e^{\mu _A}(i1),`$ (8)
$`{\displaystyle \frac{X_{(0,j)}}{X_B}}`$ $`=`$ $`{\displaystyle \frac{\mu _B^{j1}}{(j1)!}}e^{\mu _B}(j1)`$ (9)
where $`\mu _A\delta _At`$, $`\mu _B\delta _Bt`$ are mean values of each Poisson distribution at time $`t`$ and $`\delta _A`$, $`\delta _B`$ rate constants of additions of each sand. Rate equations for the reduced variables read
$`{\displaystyle \frac{dX_0}{dt}}`$ $`=`$ $`2X_0(Z_0+Z_A+Z_B)+Z_0`$ (10)
$`(\delta _A+\delta _B)X_0,`$
$`{\displaystyle \frac{dX_A}{dt}}`$ $`=`$ $`(X_0X_A)(Z_0+Z_A+Z_B)+Z_A`$ (11)
$`+\delta _AX_0\delta _BX_A,`$
$`{\displaystyle \frac{dX_B}{dt}}`$ $`=`$ $`(X_0X_B)(Z_0+Z_A+Z_B)+Z_B`$ (12)
$`+\delta _BX_0\delta _AX_B,`$
$`{\displaystyle \frac{dZ_0}{dt}}`$ $`=`$ $`(\alpha _AX_A+\alpha _BX_B2Z_0)(Z_0+Z_A+Z_B)Z_0`$ (13)
$`+\alpha _B\delta _AX_B+\alpha _A\delta _BX_A,`$
$`{\displaystyle \frac{dZ_A}{dt}}`$ $`=`$ $`\{(1\alpha _A)X_A+Z_0\}(Z_0+Z_A+Z_B)Z_A`$ (14)
$`+(1\alpha _A)\delta _BX_A,`$
$`{\displaystyle \frac{dZ_B}{dt}}`$ $`=`$ $`\{(1\alpha _B)X_B+Z_0\}(Z_0+Z_A+Z_B)Z_B`$ (15)
$`+(1\alpha _B)\delta _AX_B`$
where $`\alpha _A(t)X_{(1,0)}/X_A=\mathrm{exp}(\delta _At)`$ and $`\alpha _B(t)X_{(0,1)}/X_B=\mathrm{exp}(\delta _Bt)`$.
Numerical solutions of the equation system (10)-(15) are plotted in Fig. 3.
It is found that in the long time limit, the system goes to a steady state which is supposed to be SOC. The steady state is able to be examined analytically. When $`t\mathrm{}`$, then $`\alpha _A,\alpha _B0`$ and $`X_0,Z_00`$. Therefore the following relations are derived from Eqs. (10)-(15).
$`X_A={\displaystyle \frac{Z_A}{Z_A+Z_B+\delta _B}},X_B={\displaystyle \frac{Z_B}{Z_A+Z_B+\delta _A}}.`$ (16)
Suppose that $`Z_A+Z_B\delta _A,\delta _B`$ and $`Z_A:Z_B=\delta _A:\delta _B`$, we obtain
$`Z_A={\displaystyle \frac{\sqrt{2}\delta _B^{\frac{1}{2}}\delta _A^{\frac{3}{2}}}{(\delta _A+\delta _B)^{\frac{3}{2}}}},Z_B={\displaystyle \frac{\sqrt{2}\delta _A^{\frac{1}{2}}\delta _B^{\frac{3}{2}}}{(\delta _A+\delta _B)^{\frac{3}{2}}}}.`$ (17)
Equations (17) show that the model goes to a critical state $`(Z_A=Z_B=0)`$ in the limit $`\delta _A,\delta _B\delta 0`$. Out of two conditions (2) for SOC of the first kind, our model can successfully remove one of them with respect to the control parameter of dissipation $`ϵ`$. Furthermore, the order parameters behave as $`Z_A,Z_B\delta ^{1/2}`$ when the control parameters go to zero $`\delta 0`$. These two characteristic features of two-component sandpile models, that is the existence of the steady state without dissipation and the asymptotic behavior to the critical state in $`\delta 0`$, distinctively differ from usual SOC models and are strong evidences that the model belongs to different universality class. Consequently, we are successfully able to find a new mechanism and class of SOC. So, we call this SOC of the second kind.
We have investigated two-component sandpile models with the rule (b) and periodic boundary conditions mainly on $`1`$dimensional lattices. Extention to two components is essential to generate an infinite number of stable states which substitute for boundary dissipation. Therefore, the model becomes more natural because the dissipation is usually introduced artificially only for stopping avahanches. Simulation results indicate the model shows SOC behavior. The mean-field theory confirms that the model goes to the critical steady state belonging to a different universality class from one-component models. Here, only rough-tuning of rate constants of additions $`\delta `$ is enough to induce SOC. We can successfully construct a new SOC model without introducing any artificial dissipation mechanisms. Consequently, the two conditions (2) are not the necessary condition for SOC, which is the answer to the question mentioned previously.
It is interesting to summarize SOC models from the standpoint of the number of conserved quantities. Models that have no conserved quantities, such as contact processes, show power laws strictly on the critical point. Therefore, fine-tuning of control parameters is required. Models that have one conserved quantity, such as one-component sandpile models, become SOC in rough-tuning of two control parameters, an addition and a dissipation rates of sands, around zero. Two-component sandpile models with two conserved quantities show SOC in rough-tuning of only one control parameter, an addition rate of sands, around zero. It could be concluded that the more conservation laws result in the less control parameters.
Sandpile models are considered to have some relevance to power laws of earthequake magnitude distributions, called Guthenberg-Richter lawTurcotte1999 . One-component sandpile models have to introduce some artificial boundary dissipations in order to stop avalanches. However, the surface of the earth has no apparent boundary. At this point, it is more natural to apply two-component sandpile models to earthequake magnitude distributions. Furthermore, an earthquake could be triggered by multiple physical quantities, such as elastic energy and stress of earth’s crust. If an earthquake takes place only when both the energy and stress reach threshold simultaneously, the distributions of earthquakes would be well described by the SOC of the second kind. In this letter, we express the process in terms of a sandpile. Obviously, sands are merely symbols and could be any kind of substances which trigger the dynamics. In addition to earthquakes, therefore, a wide variety of SOC phenomena triggered by multiple factors would belong to SOC of the second kind. Future studies will make clear these points. |
warning/0506/hep-th0506086.html | ar5iv | text | # Massive Relativistic Particle Models with Bosonic Counterpart of Supersymmetry
## 1 Introduction
Important extensions of the relativistic symmetries were considered in the following two directions:
1. Supersymmetric extension, relating by supersymmetry (SUSY) transformations integer and half–integer spin fields (see e. g. ). The geometric way of describing supersymmetric multiplets is realized in terms of superfields – the functions on superspace $`Y_A=(x_\mu ,\theta _\alpha ^i,\overline{\theta }_{\dot{\alpha }}^i)`$ where $`\theta _\alpha ^i`$ are anticommuting Grassmann spinors.
2. Introduction of higher spin (HS) algebras, which act on infinite spin multiplets or if $`m=0`$ on infinite helicity multiplets (see e. g. -). The representation spaces of HS algebras are described by the functions on ‘bosonic’ superspace $`Z_A=(x_\mu ,\lambda _\alpha ^i,\overline{\lambda }_{\dot{\alpha }}^i)`$ with additional commuting spinor variables $`\lambda _\alpha ^i`$. The bosonic counterparts of superfields one can call the spinorial Kaluza–Klein (KK) fields, with spinorial additional dimensions. We obtain
$$\mathrm{\Phi }_A(Z_A)=\underset{n,k=0}{\overset{\mathrm{}}{}}\underset{_{(\dot{\beta }_1\mathrm{}\dot{\beta }_k)}^{(\alpha _1\mathrm{}\alpha _n)}}{}\underset{_{(j_1\mathrm{}j_k)}^{(i_1\mathrm{}i_n)}}{}\phi _{A;i_1\mathrm{}i_nj_1\mathrm{}j_k}^{\alpha _1\mathrm{}\alpha _n\dot{\beta }_1\mathrm{}\dot{\beta }_k}(x)\lambda _{\alpha _1}^{i_1}\mathrm{}\lambda _{\alpha _n}^{i_n}\overline{\lambda }_{\dot{\beta }_1}^{j_1}\mathrm{}\overline{\lambda }_{\dot{\beta }_k}^{j_k}.$$
(1)
The auxiliary commuting spinorial variables $`(\lambda _\alpha ^i,\overline{\lambda }_{\dot{\alpha }}^i)`$ occurs in several geometric frameworks, for example in twistor approach to the space–time geometry - or in the models with double (target and world volume) supersymmetry -.
In this note we would like to study the group-theoretic and dynamical consequences of introducing bosonic counterpart of supersymmetry, obtained by supplementing the Poincaré algebra by bosonic spinorial charges. We recall the general $`N=1`$ SUSY relation with tensorial charges
$$\{Q_a,Q_b\}=2(\gamma ^\mu C)_{ab}P_\mu +(\sigma ^{\mu \nu }C)_{ab}Z_{\mu \nu }$$
(2)
where in Majorana representation $`C=\gamma _0`$ and
$`Q_a`$ is a four–component Majorana spinor of supercharges
$`Z_{\mu \nu }=Z_{\nu \mu }`$ describe six Abelian tensorial charges.
The bosonic counterpart of general $`N=1`$ SUSY takes the form
$$[R_a,R_b]=2(\gamma ^\mu \gamma _5C)_{ab}P_\mu +2C_{ab}Z^{(1)}+2(\gamma _5C)_{ab}Z^{(2)}$$
(3)
where
$`R_a`$ is a four–component spinor of bosonic charges
$`Z^{(1)}`$ ($`Z^{(2)}`$) are scalar (pseudoscalar) central charges.
In order to obtain in (3) the standard inversion properties of the fourmomentum generator one should assume suitable transformation properties of the spinor $`R_a`$.<sup>1</sup><sup>1</sup>1Spinorial supercharges transform under space–time inversions in standard way ($`Q_a^{}=(\gamma _0Q)_a`$ for the space inversion $`P`$, $`Q_a^{}=(\gamma _0\gamma _5Q)_a`$ for the time inversion $`T`$). The bosonic spinorial charges $`R_a`$ are so–called pseudospinor ) transforming under inversion in alternative way ($`R_a^{}=(\gamma _0\gamma _5R)_a`$ under $`P`$, $`R_a^{}=(\gamma _0R)_a`$ under $`T`$).
Our aim is to study the massive relativistic particle models invariant under bosonic counterpart of SUSY and perform their quantization. Contrary to the case of simple SUSY the $`N=1`$ relation (3) contains scalar and pseudoscalar central charges, which can be related with the mass parameter. In Sect. 2 we describe (using two–component Weyl notation) the particle model describing the trajectory in the spinorial KK space $`^{4,4}`$ with the coordinates $`Z_A=(x_\mu ,\lambda _\alpha ,\overline{\lambda }_{\dot{\alpha }})`$. After calculating the complete set of constraints we perform the quantization using either Heisenberg picture or the Gupta–Bleuler method (Schrödinger picture).<sup>2</sup><sup>2</sup>2Gupta–Bleuler method has been applied to massive relativistic superparticle e. g. in . We shall obtain the wave function $`\mathrm{\Psi }(Z_A)`$ satisfying the KG equation and the bosonic counterpart of the chirality condition. In Sect. 3 we analyze the massless limit of our model, with massless fields with arbitrary helicity satisfying Fierz-Pauli equations. In Sect. 4 we consider the relativistic particle in $`N=2`$ spinorial KK space $`^{4,8}`$ with the coordinates $`(x_\mu ,\lambda _{\alpha i},\overline{\lambda }_{\dot{\alpha }i})`$ $`(i=1,2)`$. It appears that for the particular choice of bosonic counterpart of $`N=2`$ SUSY, with internal symmetry $`O(1,1)`$, one can obtain the linear Bargman–Wigner equations for $`D=4`$ massive higher spin fields . In Sect. 5 we shall discuss the problem of nonstandard relation between spin and statistics for the field components of spinorial KK fields.
## 2 Massive particle model with $`N=1`$ bosonic counterpart of SUSY.
### 2.1 Classical model
We consider the following action<sup>3</sup><sup>3</sup>3We use following notations. The metric has mostly minus $`\eta _{\mu \nu }=\mathrm{diag}(+)`$. The Weyl two–spinor indices are risen and lowered by $`\phi ^\alpha =ϵ^{\alpha \beta }\phi _\beta `$, $`\phi _\alpha =\phi ^\beta ϵ_{\beta \alpha }`$, $`\overline{\phi }^{\dot{\alpha }}=ϵ^{\dot{\alpha }\dot{\beta }}\overline{\phi }_{\dot{\beta }}`$, $`\overline{\phi }_{\dot{\alpha }}=\overline{\phi }^{\dot{\beta }}ϵ_{\dot{\beta }\dot{\alpha }}`$ where $`ϵ^{\alpha \beta }ϵ_{\beta \gamma }=\delta _\gamma ^\alpha `$, $`ϵ^{\dot{\alpha }\dot{\beta }}ϵ_{\dot{\beta }\dot{\gamma }}=\delta _{\dot{\gamma }}^{\dot{\alpha }}`$. Algebra $`\sigma `$–matrices $`\sigma _{\alpha \dot{\beta }}^\mu =(\overline{\sigma _{\beta \dot{\alpha }}^\mu })`$ and $`\sigma _\mu ^{\dot{\alpha }\alpha }=ϵ^{\alpha \beta }ϵ^{\dot{\alpha }\dot{\beta }}\sigma _{\mu \beta \dot{\beta }}`$ is $`\sigma _{\mu \alpha \dot{\gamma }}\sigma _\nu ^{\dot{\gamma }\beta }+\sigma _{\nu \alpha \dot{\gamma }}\sigma _\mu ^{\dot{\gamma }\beta }=2\eta _{\mu \nu }\delta _\alpha ^\beta `$. Also we define $`p_{\alpha \dot{\beta }}=p_\mu \sigma _{\alpha \dot{\beta }}^\mu `$, $`p^{\dot{\alpha }\beta }=p^\mu \sigma _\mu ^{\dot{\alpha }\beta }`$ for any vector $`p_\mu `$.
$$S=𝑑\tau ,$$
(4)
$$=m(\dot{\omega }_\mu \dot{\omega }^\mu )^{1/2}i(z\dot{\lambda }^\alpha \lambda _\alpha \overline{z}\overline{\lambda }_{\dot{\alpha }}\dot{\overline{\lambda }}{}_{}{}^{\dot{\alpha }})$$
(5)
where
$$d\omega ^\mu =\dot{\omega }^\mu d\tau =dx^\mu id\lambda ^\alpha \sigma _{\alpha \dot{\beta }}^\mu \overline{\lambda }^{\dot{\beta }}+i\lambda ^\alpha \sigma _{\alpha \dot{\beta }}^\mu d\overline{\lambda }^{\dot{\beta }}.$$
(6)
The action (4)-(6) describes the particle trajectory in Minkowski space extended by two commuting complex Weyl spinor coordinates $`\lambda ^\alpha (\tau )`$, $`\overline{\lambda }^{\dot{\alpha }}=(\overline{\lambda ^\alpha })`$ and invariant under the following spinorial bosonic transformation
$$\delta x^\mu =i\lambda ^\alpha \sigma _{\alpha \dot{\beta }}^\mu \overline{\epsilon }^{\dot{\beta }}i\epsilon ^\alpha \sigma _{\alpha \dot{\beta }}^\mu \overline{\lambda }^{\dot{\beta }},\delta \lambda ^\alpha =\epsilon ^\alpha ,\delta \overline{\lambda }^{\dot{\alpha }}=\overline{\epsilon }^{\dot{\alpha }}$$
(7)
where $`\epsilon ^\alpha `$ is a constant commuting Weyl spinor. The constant $`m`$ is the mass of particle whereas $`z`$ is an arbitrary complex parameter with the dimension of mass. It is easy to see that performing the suitable phase transformation $`\lambda _\alpha ^{}=e^{ia}\lambda _\alpha `$, $`\overline{\lambda }_{\dot{\alpha }}^{}=e^{ia}\overline{\lambda }_{\dot{\alpha }}`$, where $`a=\frac{1}{2}argz`$ one gets the real parameter $`z`$.
Conserved Noether spinorial charges corresponding to the transformations (7) are
$$R_\alpha \pi _\alpha ip_{\alpha \dot{\beta }}\overline{\lambda }^{\dot{\beta }}iz\lambda _\alpha ,$$
(8)
$$\overline{R}_{\dot{\alpha }}\overline{\pi }_{\dot{\alpha }}+i\lambda ^\beta p_{\beta \dot{\alpha }}+i\overline{z}\overline{\lambda }_{\dot{\alpha }}$$
(9)
where the canonical momenta are defined by
$$p_\mu =\frac{}{\dot{x}^\mu }=m(\dot{\omega }_\nu \dot{\omega }^\nu )^{1/2}\dot{\omega }_\mu ,$$
(10)
$$\pi _\alpha =\frac{}{\dot{\lambda }^\alpha }=ip_{\alpha \dot{\beta }}\overline{\lambda }^{\dot{\beta }}iz\lambda _\alpha ,$$
(11)
$$\overline{\pi }_{\dot{\alpha }}=\frac{}{\dot{\overline{\lambda }}^{\dot{\alpha }}}=i\lambda ^\beta p_{\beta \dot{\alpha }}+i\overline{z}\overline{\lambda }_{\dot{\alpha }}.$$
(12)
Using the canonical Poisson brackets
$$\{x^\mu ,p_\nu \}=\delta _\nu ^\mu ,\{\lambda ^\alpha ,\pi _\beta \}=\delta _\beta ^\alpha ,\{\overline{\lambda }^{\dot{\alpha }},\overline{\pi }_{\dot{\beta }}\}=\delta _{\dot{\beta }}^{\dot{\alpha }}$$
(13)
we obtain the algebra
$$\{R_\alpha ,\overline{R}_{\dot{\beta }}\}=2ip_{\alpha \dot{\beta }},$$
(14)
$$\{R_\alpha ,R_\beta \}=2izϵ_{\alpha \beta },\{\overline{R}_{\dot{\alpha }},\overline{R}_{\dot{\beta }}\}=2i\overline{z}ϵ_{\dot{\alpha }\dot{\beta }}$$
(15)
which is equivalent to the algebra (3) with $`Z=Z^{(1)}+iZ^{(2)}=z`$.
From (10)–(12) follow the mass shell constraint and the set of four spinorial constraints
$$Tp^2m^20,$$
(16)
$$D_\alpha \pi _\alpha +ip_{\alpha \dot{\beta }}\overline{\lambda }^{\dot{\beta }}+iz\lambda _\alpha 0,$$
(17)
$$\overline{D}_{\dot{\alpha }}\overline{\pi }_{\dot{\alpha }}i\lambda ^\beta p_{\beta \dot{\alpha }}i\overline{z}\overline{\lambda }_{\dot{\alpha }}0$$
(18)
Using the formulae (10)-(12) we confirm that the canonical Hamiltonian vanishes<sup>4</sup><sup>4</sup>4The vanishing of Hamiltonian follows from the invariance of the action (46) under the arbitrary local rescaling $`\tau \tau ^{}=\tau ^{}(\tau )`$.
$$=\dot{x}^\mu p_\mu +\dot{\lambda }^\alpha \pi _\alpha +\overline{\pi }_{\dot{\alpha }}\dot{\overline{\lambda }}{}_{}{}^{\dot{\alpha }}=0.$$
and the total Hamiltonian is the linear combination of first class constraints multiplied by Lagrange multipliers.
The constraints (16)-(18) satisfy the following Poisson brackets
$$\{D_\alpha ,\overline{D}_{\dot{\beta }}\}=2ip_{\alpha \dot{\beta }},$$
(19)
$$\{D_\alpha ,D_\beta \}=2izϵ_{\alpha \beta },\{\overline{D}_{\dot{\alpha }},\overline{D}_{\dot{\beta }}\}=2i\overline{z}ϵ_{\dot{\alpha }\dot{\beta }}.$$
(20)
The scalar constraint $`T0`$ is first class and all the spinorial constraints (17), (18) are second class. Indeed we find that the determinant of the matrix
$$𝒞=\left(\begin{array}{cc}\{D_\alpha ,D_\beta \}& \{D_\alpha ,\overline{D}_{\dot{\beta }}\}\\ \{\overline{D}_{\dot{\alpha }},D_\beta \}& \{\overline{D}_{\dot{\alpha }},\overline{D}_{\dot{\beta }}\}\end{array}\right)=\left(\begin{array}{cc}2izϵ_{\alpha \beta }& 2ip_{\alpha \dot{\beta }}\\ 2ip_{\beta \dot{\alpha }}& 2i\overline{z}ϵ_{\dot{\alpha }\dot{\beta }}\end{array}\right).$$
(21)
is equal to<sup>5</sup><sup>5</sup>5 In calculation it is convenient to use that $`det𝒞=detDdet(ABD^1C)`$ for matrix $`𝒞=\left(\begin{array}{cc}A& B\\ C& D\end{array}\right)=\left(\begin{array}{cc}1& B\\ 0& D\end{array}\right)\left(\begin{array}{cc}ABD^1C& 0\\ D^1C& 1\end{array}\right)`$.
$$det𝒞=16(p^2+|z|^2)^2.$$
(22)
We see from (22) that the matrix $`𝒞`$ (see (21)) is invertible for any $`z`$, and the constraints (1718) are second class.
### 2.2 Quantization
The first quantization of the model can be performed using one of two methods:
i) Following the technique of quantization of systems with second class constraints one can introduce Dirac brackets (DB) for the independent phase space degrees of freedom $`𝒵_M=(x^\mu ,p_\mu ,\lambda _\alpha ,\overline{\lambda }_{\dot{\alpha }})`$
$$\{𝒵_M,𝒵_N\}^{}=\{𝒵_M,𝒵_N\}\{𝒵_M,D_r\}(𝒞^1)_{rs}\{D_s,𝒵_N\}$$
(23)
where $`D_r=(D_\alpha ,\overline{D}_{\dot{\alpha }})`$. In particular for suitably normalized spinor coordinates<sup>6</sup><sup>6</sup>6On the mass shell $`T0`$ and at $`z=m`$ we get $`\eta _\alpha =2m\lambda _\alpha `$ and $`\overline{\eta }_{\dot{\alpha }}=2m\overline{\lambda }_{\dot{\alpha }}`$.
$$\eta _\alpha =[2(p^2+|z|^2)]^{1/2}\lambda _\alpha ,\overline{\eta }_{\dot{\alpha }}=[2(p^2+|z|^2)]^{1/2}\overline{\lambda }_{\dot{\alpha }}$$
(24)
one obtains the relations
$$\{\eta _\alpha ,\eta _\beta \}^{}=i\overline{z}ϵ_{\alpha \beta },\{\overline{\eta }_{\dot{\alpha }},\overline{\eta }_{\dot{\beta }}\}^{}=izϵ_{\dot{\alpha }\dot{\beta }},\{\eta _\alpha ,\overline{\eta }_{\dot{\beta }}\}^{}=ip_{\alpha \dot{\beta }}$$
(25)
leading after quantization to noncommutative Weyl spinor coordinates. Similarly one can calculate
$$\{x_\mu ,x_\nu \}^{}=\frac{i}{2(p^2+|z|^2)}S_{\mu \nu },$$
$$S^{\mu \nu }=\lambda ^\alpha [(\sigma ^{\mu \nu })_\alpha {}_{}{}^{\beta }p_{\beta \dot{\gamma }}^{}+p_{\alpha \dot{\beta }}(\overline{\sigma }^{\mu \nu })^{\dot{\beta }}{}_{\dot{\gamma }}{}^{}]\overline{\lambda }^{\dot{\gamma }}+z\lambda ^\alpha (\sigma ^{\mu \nu })_\alpha {}_{}{}^{\beta }\lambda _{\beta }^{}+\overline{z}\overline{\lambda }_{\dot{\alpha }}(\overline{\sigma }^{\mu \nu })^{\dot{\alpha }}{}_{\dot{\beta }}{}^{}\overline{\lambda }_{}^{\dot{\beta }},$$
$$(\sigma ^{\mu \nu })_\alpha {}_{}{}^{\beta }\frac{1}{2}(\sigma _{\alpha \dot{\gamma }}^\mu \sigma ^{\nu \dot{\gamma }\beta }\sigma _{\alpha \dot{\gamma }}^\nu \sigma ^{\mu \dot{\gamma }\beta }),(\overline{\sigma }^{\mu \nu })^{\dot{\alpha }}{}_{\dot{\beta }}{}^{}\frac{1}{2}(\sigma ^{\mu \dot{\alpha }\gamma }\sigma _{\gamma \dot{\beta }}^\nu \sigma ^{\nu \dot{\alpha }\gamma }\sigma _{\gamma \dot{\beta }}^\mu ),$$
i.e. we see that the coordinates are becoming also noncommutative.
One can note that after the linear transformation of the form
$$\eta _\alpha ^{}=\eta _\alpha +cp_{\alpha \dot{\beta }}\overline{\eta }^{\dot{\beta }},\overline{\eta }_\alpha ^{}=\overline{\eta }_{\dot{\alpha }}+\overline{c}\eta ^\beta p_{\beta \dot{\alpha }},$$
(26)
we can obtain from the algebra (25) for certain choice of $`c`$ the DB relations $`\{\eta _\alpha ^{},\eta _\beta ^{}\}^{}ϵ_{\alpha \beta }`$, $`\{\eta _\alpha ^{},\overline{\eta }_{\dot{\beta }}^{}\}^{}=0`$. The algebra of such type is used for description of massless fields with arbitrary helicities in . For other choice of $`c`$ we obtain alternatively $`\{\eta _\alpha ^{},\eta _\beta ^{}\}^{}=0`$, $`\{\eta _\alpha ^{},\overline{\eta }_{\dot{\beta }}^{}\}^{}p_{\alpha \dot{\beta }}`$. In such a case $`\eta _\alpha ^{}`$ and $`\overline{\eta }_{\dot{\alpha }}^{}`$ can be treated as of suitably rescaled creation and annihilation operators.
ii) Other way is the Gupta–Bleuler quantization method. Such a technique implies the split of the second class constraints into complex–conjugated pairs, with holomorphic and antiholomorphic parts forming separately the subalgebras of first class constraints. The algebra (1920) of the constraints (17), (18) does not satisfy these requirements. Let us introduce, however, new constraints as follows
$$𝒟_\alpha =D_\alpha +\frac{b}{\overline{z}}p_{\alpha \dot{\beta }}\overline{D}^{\dot{\beta }},\overline{𝒟}_{\dot{\alpha }}=\overline{D}_{\dot{\alpha }}+\frac{\overline{b}}{z}D^\beta p_{\beta \dot{\alpha }},$$
(27)
$`\overline{𝒟}_{\dot{\alpha }}=(\overline{𝒟_\alpha })`$. If $`b`$ satisfies the equation $`(b^22b)\frac{m^2}{|z|^2}1=0`$ (i.e. $`b=(1\pm \sqrt{1+\frac{|z|^2}{m^2}})`$) the algebra of the constraints (27) takes the form
$$\{𝒟_\alpha ,𝒟_\beta \}=\frac{2i}{\overline{z}}ϵ_{\alpha \beta }T,\{\overline{𝒟}_{\dot{\alpha }},\overline{𝒟}_{\dot{\beta }}\}=\frac{2i}{z}ϵ_{\dot{\alpha }\dot{\beta }}T,\{𝒟_\alpha ,\overline{𝒟}_{\dot{\beta }}\}=4b(1+\frac{m^2}{|z|^2})ip_{\alpha \dot{\beta }}\frac{2b^2i}{|z|^2}p_{\alpha \dot{\beta }}T.$$
We see that the constraints (27) are suitable for application of Gupta–Bleuler quantization method. It should be mentioned that the transformation from constraints $`(D_\alpha ,\overline{D}_{\dot{\alpha }})`$ to constraints $`(𝒟_\alpha ,\overline{𝒟}_{\dot{\alpha }})`$ is invertible.
We shall assume that the wave function satisfies the Klein–Gordon equation, what follows from the constraint (16). On the mass shell (16) the constraints (27) have the form
$$𝒟_\alpha =\pi _\alpha ^{}2b(1+\frac{m^2}{|z|^2})ip_{\alpha \dot{\beta }}\overline{\lambda }^{\dot{\beta }}0,\overline{𝒟}_{\dot{\alpha }}=\overline{\pi }_{\dot{\alpha }}^{}+2b(1+\frac{m^2}{|z|^2})i\lambda ^\beta p_{\beta \dot{\alpha }}0$$
(28)
where we introduced new spinor variables via the following canonical transformation
$$\pi _\alpha ^{}\pi _\alpha +\frac{b}{\overline{z}}p_{\alpha \dot{\beta }}\overline{\pi }^{\dot{\beta }},\overline{\pi }_{\dot{\alpha }}^{}\overline{\pi }_{\dot{\alpha }}+\frac{b}{z}\pi ^\beta p_{\beta \dot{\alpha }},$$
(29)
$$\lambda ^\alpha \frac{|z|^2}{|z|^2+b^2p^2}(\lambda ^\alpha \frac{b}{z}\overline{\lambda }_{\dot{\beta }}p^{\dot{\beta }\alpha }),\overline{\lambda }^{\dot{\alpha }}\frac{|z|^2}{|z|^2+b^2p^2}(\overline{\lambda }^{\dot{\alpha }}\frac{b}{\overline{z}}p^{\dot{\alpha }\beta }\lambda _\beta )$$
(30)
i.e. we obtain the standard canonical commutation relations (compare with (13))
$$\{\lambda ^\alpha ,\pi _\beta ^{}\}=\delta _\beta ^\alpha ,\{\overline{\lambda }^{\dot{\alpha }},\overline{\pi }_{\dot{\beta }}^{}\}=\delta _{\dot{\beta }}^{\dot{\alpha }},\{\lambda ^\alpha ,\overline{\pi }_{\dot{\beta }}^{}\}=\{\overline{\lambda }^{\dot{\alpha }},\pi _\beta ^{}\}=0.$$
(31)
For the quantization of our model we consider the Schrödinger representation of the CCR (31)
$$\pi _\alpha ^{}=i/\lambda ^\alpha ,\overline{\pi }_{\dot{\alpha }}^{}=i/\overline{\lambda }^{\dot{\alpha }}$$
(32)
and use the wave function $`\mathrm{\Psi }`$ in the momentum representation, i.e. $`\mathrm{\Psi }=\mathrm{\Psi }(p_\mu ,\lambda ^\alpha ,\overline{\lambda }^{\dot{\alpha }})`$. The spinorial wave equation $`\overline{𝒟}_{\dot{\alpha }}\mathrm{\Psi }=0`$ takes the following form<sup>7</sup><sup>7</sup>7The choice of $`𝒟_\alpha `$ in place $`\overline{𝒟}_{\dot{\alpha }}`$ is equally well possible.
$$(/\overline{\lambda }^{\dot{\alpha }}+2b(1+\frac{m^2}{|z|^2})\lambda ^\beta p_{\beta \dot{\alpha }})\mathrm{\Psi }=0.$$
(33)
The solution of (33) is given by
$$\mathrm{\Psi }(p_\mu ,\lambda ^\alpha ,\overline{\lambda }^{\dot{\alpha }})=e^{2b(1+{\scriptscriptstyle \frac{m^2}{|z|^2}})\lambda ^\beta p_{\beta \dot{\alpha }}\overline{\lambda }^{\dot{\alpha }}}\stackrel{~}{\mathrm{\Psi }}(p_\mu ,\lambda ^\alpha )$$
(34)
where the field $`\stackrel{~}{\mathrm{\Psi }}(p_\mu ,\lambda ^\alpha )`$ depends only on one Weyl spinor $`\lambda ^\alpha `$ and provides the bosonic counterpart of $`D=4`$ $`N=1`$ chiral superfield.
Due to the bosonic nature of $`\lambda ^\alpha `$ in expansion of $`\stackrel{~}{\mathrm{\Psi }}(p_\mu ,\lambda ^\alpha )`$ there is an infinite number of space–time fields $`\psi _{\alpha _1\mathrm{}\alpha _n}(p)=\psi _{(\alpha _1\mathrm{}\alpha _n)}(p)`$, $`n=0,1,\mathrm{},\mathrm{}`$. The mass–shell condition (16) after the transition by Fourier transformation to the space–time picture, leads to the Klein–Gordon (KG) equation ($`\mathrm{}_\mu ^\mu `$)
$$(\mathrm{}+m^2)\mathrm{\Psi }(x;\lambda ,\overline{\lambda })=0(\mathrm{}+m^2)\psi _{\alpha _1\mathrm{}\alpha _n}(x)=0(n=0,1,2,\mathrm{}).$$
(35)
Here we should observe that
i) The half-integer spin fields (n odd) satisfy KG equation, however in massless case the half-integer helicity fields do satisfy linear equations (see Sect. 3),
ii) The spin–statistic theorem is not valid – both integer and half–integer spin fields are bosonic. We shall come back to the question of statistics in Sect. 4.
## 3 Massless particle model with $`N=1`$ bosonic counterpart of SUSY
The model (4), (5) can be described equivalently by the Lagrangian
$$=\frac{1}{2e}(\dot{\omega }_\mu \dot{\omega }^\mu +e^2m^2)i(z\dot{\lambda }^\alpha \lambda _\alpha \overline{z}\overline{\lambda }_{\dot{\alpha }}\dot{\overline{\lambda }}{}_{}{}^{\dot{\alpha }}).$$
(36)
After eliminating the einbein $`e`$ by means of its equation of motion, from (36) one obtains the Lagrangian (5). The massless limit of (36) looks as follows
$$=\frac{1}{2e}\dot{\omega }_\mu \dot{\omega }^\mu i(z\dot{\lambda }^\alpha \lambda _\alpha \overline{z}\overline{\lambda }_{\dot{\alpha }}\dot{\overline{\lambda }}{}_{}{}^{\dot{\alpha }}).$$
(37)
Besides the constraint $`p_e0`$ which implies pure gauge character of the einbein $`e`$, from (37) one gets the following constraints
$$T=p^20,$$
(38)
$$D_\alpha =\pi _\alpha +ip_{\alpha \dot{\beta }}\overline{\lambda }^{\dot{\beta }}+iz\lambda _\alpha 0,\overline{D}_{\dot{\alpha }}=\overline{\pi }_{\dot{\alpha }}i\lambda ^\beta p_{\beta \dot{\alpha }}i\overline{z}\overline{\lambda }_{\dot{\alpha }}0.$$
(39)
The nonvanishing Poisson brackets are
$$\{D_\alpha ,\overline{D}_{\dot{\beta }}\}=2ip_{\alpha \dot{\beta }},\{D_\alpha ,D_\beta \}=2izϵ_{\alpha \beta },\{\overline{D}_{\dot{\alpha }},\overline{D}_{\dot{\beta }}\}=2i\overline{z}ϵ_{\dot{\alpha }\dot{\beta }}.$$
(40)
The mass constraint (38) is of the first class. The determinant of the Poisson brackets matrix (21) characterizing the spinorial constraints (39) is the following
$$det𝒞=16(p^2+|z|^2)^216|z|^4.$$
(41)
If $`z0`$ all spinorial constraints (39) are second class. In the case of vanishing central charges $`z=0`$ the four spinorial constraints (39) contain two second class constraints and two first class. Below we analyze massless particle at $`z=0`$ with spinorial first class constraints, defined as follows:
$$F^{\dot{\alpha }}=p^{\dot{\alpha }\beta }D_\beta 0,\overline{F}^\alpha =\overline{D}_{\dot{\beta }}p^{\dot{\beta }\alpha }0$$
(42)
with the following Poisson brackets
$$\{F^{\dot{\alpha }},\overline{D}_{\dot{\beta }}\}=2i\delta _{\dot{\beta }}^{\dot{\alpha }}T0,\{\overline{F}^\alpha ,D_\beta \}=2i\delta _\beta ^\alpha T0$$
But the first class constraints (42) are reducible: $`p_{\alpha \dot{\beta }}F^{\dot{\beta }}0`$, $`\overline{F}^\beta p_{\beta \dot{\alpha }}0`$. The irreducible separation of first and second class constraints is obtained by projecting of the spinorial constraints (39) along spinors $`\lambda ^\alpha `$ and $`\overline{\lambda }_{\dot{\alpha }}p^{\dot{\alpha }\alpha }`$.<sup>8</sup><sup>8</sup>8This procedure is corrected since spinors $`\lambda ^\alpha `$ and $`\overline{\lambda }_{\dot{\alpha }}p^{\dot{\alpha }\alpha }`$ are not proportional in considered task. Otherwise, when $`\lambda ^\alpha p_{\alpha \dot{\alpha }}\overline{\lambda }^{\dot{\alpha }}=0`$, we have $`p_{\alpha \dot{\alpha }}\lambda _\alpha \overline{\lambda }_{\dot{\alpha }}`$. Then the spinorial constraints (39), taking the form $`\pi _\alpha 0`$, $`\overline{\pi }_{\dot{\alpha }}0`$, exclude completely the dependence on $`\lambda `$, $`\overline{\lambda }`$. As result we obtain the system describing only by the variables $`x^\mu `$, $`p_\mu `$ and the constraint (38) i. e. the massless particle of zero helicity. The constraints
$$G\lambda ^\alpha D_\alpha 0,\overline{G}\overline{D}_{\dot{\alpha }}\overline{\lambda }^{\dot{\alpha }}0$$
(43)
are second class whereas the constraints
$$F\overline{\lambda }_{\dot{\alpha }}p^{\dot{\alpha }\alpha }D_\alpha 0,\overline{F}\overline{D}_{\dot{\alpha }}p^{\dot{\alpha }\alpha }\lambda _\alpha 0$$
(44)
are of first class. Their Poisson brackets look as follows:
$$\{G,\overline{G}\}=2i\lambda ^\alpha p_{\alpha \dot{\alpha }}\overline{\lambda }^{\dot{\alpha }}0,\{F,\overline{F}\}=(\lambda ^\alpha \pi _\alpha \overline{\pi }_{\dot{\alpha }}\overline{\lambda }^{\dot{\alpha }})T,$$
$$\{G,F\}=\{\overline{G},F\}=F,\{G,\overline{F}\}=\{\overline{G},\overline{F}\}=\overline{F}.$$
We carry out quantization of massless particle with $`N=1`$ bosonic counterpart of SUSY by Gupta–Bleuler method. The wave equations are imposed by the first class constraints (38), (44) $`T0`$, $`F0`$, $`\overline{F}0`$ and either $`\overline{G}0`$ or $`G0`$. But the pair of constraints $`G0`$ and $`F0`$ are equivalent to the constraints $`D_\alpha 0`$; similarly the constraints $`\overline{G}0`$ and $`\overline{F}0`$ are equivalent to the constraints $`\overline{D}_{\dot{\alpha }}0`$. Thus we have two possible quantizations:
– ‘bosonic chiral’ quantization with the wave equations
$$T|\mathrm{\Psi }=0,F|\mathrm{\Psi }=0,\overline{D}_{\dot{\alpha }}|\mathrm{\Psi }=0$$
(45)
– ‘bosonic antichiral’ quantization with wave function subjected to the conditions
$$T|\mathrm{\Psi }=0,\overline{F}|\mathrm{\Psi }=0,D_\alpha |\mathrm{\Psi }=0.$$
(46)
Let us consider the chiral case. In the representation
$$p_\mu =i/x^\mu i_\mu ,\pi _\alpha =i/\lambda ^\alpha i_\alpha ,\overline{\pi }_{\dot{\alpha }}=i/\overline{\lambda }^{\dot{\alpha }}i\overline{}_{\dot{\alpha }}$$
the wave function $`\mathrm{\Psi }(x,\lambda ,\overline{\lambda })`$ satisfies the equations
$$\mathrm{}\mathrm{\Psi }=0,$$
(47)
$$\overline{D}_{\dot{\alpha }}\mathrm{\Psi }=(i\overline{}_{\dot{\alpha }}\lambda ^\beta _{\beta \dot{\alpha }})\mathrm{\Psi }=0$$
(48)
$$i\overline{\lambda }_{\dot{\alpha }}^{\dot{\alpha }\alpha }D_\alpha \mathrm{\Psi }=\overline{\lambda }_{\dot{\alpha }}^{\dot{\alpha }\alpha }_\alpha \mathrm{\Psi }=0$$
(49)
In the variables $`x_L^\mu =x^\mu +i\lambda \sigma ^\mu \overline{\lambda }`$, $`\lambda ^\alpha `$, $`\overline{\lambda }^{\dot{\alpha }}`$ bosonic SUSY-covariant derivatives take the form
$$D_\alpha =i_\alpha +2_{L\alpha \dot{\alpha }}\overline{\lambda }^{\dot{\alpha }},\overline{D}_{\dot{\alpha }}=i\overline{}_{\dot{\alpha }}.$$
(50)
Thus due to the chirality condition (48) the wave function does not depend on $`\overline{\lambda }^{\dot{\alpha }}`$. It depends only on the left chiral variables $`z_L=(x_L^\mu ,\lambda ^\alpha )`$, and commuting spinor $`\lambda `$. One can write the following expansion
$$\mathrm{\Psi }(x_L,\lambda )=\underset{n=0}{\overset{\mathrm{}}{}}\lambda ^{\alpha _1}\mathrm{}\lambda ^{\alpha _n}\varphi _{\alpha _1\mathrm{}\alpha _n}(x_L)$$
(51)
where the multispinor fields are totally symmetric in spinor indices, i.e. $`\varphi _{\alpha _1\mathrm{}\alpha _n}=\varphi _{(\alpha _1\mathrm{}\alpha _n)}`$. The usual fields depending on real space–time coordinates $`x^\mu `$ are obtained by
$$\varphi _{\alpha _1\mathrm{}\alpha _n}(x)=e^{i\lambda \sigma ^\mu \overline{\lambda }_\mu }\varphi _{\alpha _1\mathrm{}\alpha _n}(x_L).$$
The equation (49) gives Fierz–Pauli equations for the component fields
$$^{\dot{\beta }\beta }\varphi _{\beta \alpha _2\mathrm{}\alpha _n}=0.$$
(52)
The Klein–Gordon equation $`\mathrm{}\varphi _{\alpha _1\mathrm{}\alpha _n}=0`$, resulting from (47), follows also from (52). We see therefore that the expansion of the wave function (51) describes an infinite set of massless particles with helicities $`n/2`$.
The Gupta-Bleuler quantization procedure presented here is analogous to the one used for the quantization of massless Brink-Schwarz superparticle, but due to the bosonic character of spinorial variable $`\lambda _\alpha `$ we get infinite helicity spectrum. We recall that the infinite set of integer and half–integer helicities describes also the spectrum of supersymmetric massless particles propagating in tensorial superspace .
## 4 Massive relativistic particles with $`N=2`$ bosonic counterpart of SUSY.
### 4.1 $`N=2`$ action and the constraints
Let us introduce two commuting Weyl spinors $`\lambda _i^\alpha `$, $`\overline{\lambda }_i^{\dot{\alpha }}=(\overline{\lambda _i^\alpha })`$ ($`i=1,2`$). The natural generalization of the Lagrangian (5) is
$$=m(\dot{\omega }_\mu \dot{\omega }^\mu )^{1/2}i(z_{ij}\dot{\lambda }_i^\alpha \lambda _{\alpha j}\overline{z}_{ij}\overline{\lambda }_{\dot{\alpha }j}\dot{\overline{\lambda }}{}_{i}{}^{\dot{\alpha }}).$$
(53)
Here the constant matrix $`z_{ij}`$ is symmetric, $`z_{ij}=z_{ji}`$; the last terms in (53) are total derivatives, e.g. $`z_{ij}\dot{\lambda }_i^\alpha \lambda _{\alpha j}=\frac{1}{2}(z_{ij}\lambda _i^\alpha \lambda _{\alpha j})`$ if $`z_{ij}=z_{ji}`$.
The $`\omega `$–form can be written in general case as follow
$$\dot{\omega }^\mu =\dot{x}^\mu i\kappa _{ij}(\dot{\lambda }_i^\alpha \sigma _{\alpha \dot{\beta }}^\mu \overline{\lambda }_j^{\dot{\beta }}\lambda _j^\alpha \sigma _{\alpha \dot{\beta }}^\mu \dot{\overline{\lambda }}_i^{\dot{\beta }})$$
(54)
where $`\kappa _{ij}=\kappa _{ji}`$ is the $`2\times 2`$ Hermitean metric in $`N=2`$ unitary space. If we consider possible linear definitions of spinors $`\lambda _i^\alpha `$ in $`N=2`$ internal space one can choose
$$\kappa _{ij}=\left(\begin{array}{cc}1& 0\\ 0& \kappa \end{array}\right)$$
(55)
where $`\kappa `$ is real.
From expressions for the canonical momenta
$$p_\mu =\frac{}{\dot{x}^\mu }=m(\dot{\omega }_\nu \dot{\omega }^\nu )^{1/2}\dot{\omega }_\mu ,$$
(56)
$$\pi _{\alpha i}=\frac{}{\dot{\lambda }_i^\alpha }=i\kappa _{ij}p_{\alpha \dot{\beta }}\overline{\lambda }_j^{\dot{\beta }}iz_{ij}\lambda _{\alpha j},$$
(57)
$$\overline{\pi }_{\dot{\alpha }i}=\frac{}{\dot{\overline{\lambda }}_i^{\dot{\alpha }}}=i\kappa _{ij}\lambda _j^\beta p_{\beta \dot{\alpha }}+i\overline{z}_{ij}\overline{\lambda }_{\dot{\alpha }j}$$
(58)
we obtain the following constraints
$$Tp^2m^20,$$
(59)
$$D_{\alpha i}\pi _{\alpha i}+i\kappa _{ij}p_{\alpha \dot{\beta }}\overline{\lambda }_j^{\dot{\beta }}+iz_{ij}\lambda _{\alpha j}0,$$
(60)
$$\overline{D}_{\dot{\alpha }i}\overline{\pi }_{\dot{\alpha }i}i\kappa _{ij}\lambda _j^\beta p_{\beta \dot{\alpha }}i\overline{z}_{ij}\overline{\lambda }_{\dot{\alpha }j}0.$$
(61)
Using the canonical Poisson brackets
$$\{x^\mu ,p_\nu \}=\delta _\nu ^\mu ,\{\lambda _i^\alpha ,\pi _{\beta j}\}=\delta _\beta ^\alpha \delta _{ij},\{\overline{\lambda }_i^{\dot{\alpha }},\overline{\pi }_{\dot{\beta }j}\}=\delta _{\dot{\beta }}^{\dot{\alpha }}\delta _{ij}$$
($`\{\lambda _{\alpha i},\pi _{\beta j}\}=\{\pi _{\alpha i},\lambda _{\beta j}\}=ϵ_{\alpha \beta }\delta _{ij}`$, $`\{\overline{\lambda }_{\dot{\alpha }i},\overline{\pi }_{\dot{\beta }j}\}=\{\overline{\pi }_{\dot{\alpha }i},\overline{\lambda }_{\dot{\beta }j}\}=ϵ_{\dot{\alpha }\dot{\beta }}\delta _{ij}`$) we obtain nonzero Poisson brackets of the constraints (59)-(61)
$$\{D_{\alpha i},\overline{D}_{\dot{\beta }j}\}=2i\kappa _{ij}p_{\alpha \dot{\beta }},$$
(62)
$$\{D_{\alpha i},D_{\beta j}\}=2iz_{ij}ϵ_{\alpha \beta },\{\overline{D}_{\dot{\alpha }i},\overline{D}_{\dot{\beta }j}\}=2i\overline{z}_{ij}ϵ_{\dot{\alpha }\dot{\beta }}.$$
(63)
It should be pointed out that the relations (62), (63) with changed sign on the rhs describe the bosonic counterpart of the generalized $`N=2`$ superalgebra with the Hermitean metric $`\kappa _{ij}`$ in internal $`N=2`$ space.
The constraint (59) $`T0`$ is the first class constraint. From the spinor constraints (60), (61) one gets the following $`4\times 4`$ matrix of PB
$$𝒞=\left(\begin{array}{cc}\{D_{\alpha i},D_{\beta j}\}& \{D_{\alpha i},\overline{D}_{\dot{\beta }j}\}\\ \{\overline{D}_{\dot{\alpha }i},D_{\beta j}\}& \{\overline{D}_{\dot{\alpha }i},\overline{D}_{\dot{\beta }j}\}\end{array}\right)=\left(\begin{array}{cc}2iz_{ij}ϵ_{\alpha \beta }& 2i\kappa _{ij}p_{\alpha \dot{\beta }}\\ 2i\kappa _{ij}p_{\beta \dot{\alpha }}& 2i\overline{z}_{ij}ϵ_{\dot{\alpha }\dot{\beta }}\end{array}\right).$$
(64)
We obtain that
$$det𝒞=2^8[det(\widehat{z}\widehat{\overline{z}}+p^2\widehat{\kappa }\widehat{\overline{z}}^1\widehat{\kappa }\widehat{\overline{z}})]^2$$
where ‘hats’ denote the corresponding matrices, i.e. $`\widehat{z}=(z_{ij})`$, $`\widehat{\overline{z}}=(\overline{z}_{ij})`$ and $`\widehat{\kappa }=(\kappa _{ij})`$ is given by (55). One can consider two cases:
i) If matrix $`\widehat{z}=(z_{ij})`$ is diagonal, $`(z_{ij})=\left(\begin{array}{cc}z_1& 0\\ 0& z_2\end{array}\right)`$, we obtain that $`det(\widehat{z}\widehat{\overline{z}}+p^2\widehat{\kappa }\widehat{\overline{z}}^1\widehat{\kappa }\widehat{\overline{z}})=(|z_1|^2+p^2)(|z_2|^2+p^2\kappa ^2)`$, i.e. it is always nonvanishing. We see therefore that for arbitrary values of $`\kappa `$ and $`z_1,z_2`$ all the constraints (60), (61) are second class.
ii) In case of antidiagonal matrix $`(z_{ij})=\left(\begin{array}{cc}0& z\\ z& 0\end{array}\right)`$ (we remind that matrix $`z_{ij}`$ is symmetric), we obtain that $`det(\widehat{z}\widehat{\overline{z}}+p^2\widehat{\kappa }\widehat{\overline{z}}^1\widehat{\kappa }\widehat{\overline{z}})=(|z|^2+p^2\kappa )^2`$. One gets that the matrix of Poisson brackets of the constraints (64) has vanishing determinant if $`\kappa =\frac{|z|^2}{m^2}<0`$ and we conclude that in such a case the first class constraints are present in the model. Putting $`z=m`$, i.e. $`\kappa =1`$, it is easy to check that the unitary metric tensor $`\kappa _{ij}`$ implies the invariance of the form $`\omega _\mu `$ (see (54)) under $`U(1,1)`$ symmetry. The presence of the central charge reduces however this symmetry to the invariance group $`O(1,1)=U(1,1)O(2;c)`$, and only in this case the first class constraints are present in the model (53).<sup>9</sup><sup>9</sup>9We recall that in case of standard $`N=2`$ superparticle when spinor variables are Grassmannian and the matrix $`z_{ij}`$ is skew–symmetric, $`(z_{ij})=\left(\begin{array}{cc}0& z\\ z& 0\end{array}\right)`$, the first class constraints are presented (the matrix of Poisson brackets of the constraints has vanishing determinant) if $`\kappa =\frac{|z|^2}{m^2}>0`$ and the internal $`N=2`$ symmetry in the presence of central charges $`z=m`$ is $`U(2)Sp(2;c)=SU(2)`$.
In case ii) we will consider a simple choice $`z=m`$, i.e. $`\kappa =\frac{|z|^2}{m^2}=1`$. Introducing the notations $`\lambda _1^\alpha \lambda ^\alpha `$ and $`\lambda _2^\alpha \eta ^\alpha `$ the Lagrangian (53) and $`\omega `$–form (54) are
$$=m(\dot{\omega }_\mu \dot{\omega }^\mu )^{1/2}im(\dot{\lambda }^\alpha \eta _\alpha +\dot{\eta }^\alpha \lambda _\alpha \overline{\lambda }_{\dot{\alpha }}\dot{\overline{\eta }}{}_{}{}^{\dot{\alpha }}\overline{\eta }_{\dot{\alpha }}\dot{\overline{\lambda }}{}_{}{}^{\dot{\alpha }}),$$
(65)
$$\dot{\omega }^\mu =\dot{x}^\mu i(\dot{\lambda }^\alpha \sigma _{\alpha \dot{\beta }}^\mu \overline{\lambda }^{\dot{\beta }}\lambda ^\alpha \sigma _{\alpha \dot{\beta }}^\mu \dot{\overline{\lambda }}^{\dot{\beta }})+i(\dot{\eta }^\alpha \sigma _{\alpha \dot{\beta }}^\mu \overline{\eta }^{\dot{\beta }}\eta ^\alpha \sigma _{\alpha \dot{\beta }}^\mu \dot{\overline{\eta }}^{\dot{\beta }}).$$
(66)
### 4.2 Description of the model in terms of Dirac spinors
The formulation (65) has an attractive interpretation if we pass to the commuting four–component Dirac spinor
$$\psi _a=\left(\begin{array}{c}\lambda _\alpha \\ \overline{\eta }^{\dot{\alpha }}\end{array}\right)$$
where $`a=1,2,3,4`$. The Dirac matrices $`(\gamma _\mu )_a^b`$ in Weyl representation are as follows
$$(\gamma _\mu )_a{}_{}{}^{b}=\left(\begin{array}{cc}0& \sigma _{\alpha \dot{\beta }}^\mu \\ \sigma ^{\mu \dot{\alpha }\beta }& 0\end{array}\right),\{\gamma _\mu ,\gamma _\nu \}=2\eta _{\mu \nu }$$
where $`\sigma _{\alpha \dot{\beta }}^0=\sigma ^{0\dot{\alpha }\beta }=1_2`$ and $`\sigma _{\alpha \dot{\beta }}^i=\sigma ^{i\dot{\alpha }\beta }`$ $`(i=1,2,3)`$ are the Pauli matrices. Then
$$\overline{\psi }^a=(\psi ^+\gamma _0)^a=(\eta ^\alpha ,\overline{\lambda }_{\dot{\alpha }})$$
and we obtain
$$\dot{\overline{\psi }}\psi \overline{\psi }\dot{\psi }=\dot{\lambda }^\alpha \eta _\alpha +\dot{\eta }^\alpha \lambda _\alpha \overline{\lambda }_{\dot{\alpha }}\dot{\overline{\eta }}{}_{}{}^{\dot{\alpha }}\overline{\eta }_{\dot{\alpha }}\dot{\overline{\lambda }}{}_{}{}^{\dot{\alpha }},$$
$$\dot{\overline{\psi }}\gamma ^\mu \psi \overline{\psi }\gamma ^\mu \dot{\psi }=\dot{\eta }^\alpha \sigma _{\alpha \dot{\beta }}^\mu \overline{\eta }^{\dot{\beta }}\dot{\lambda }^\alpha \sigma _{\alpha \dot{\beta }}^\mu \overline{\lambda }^{\dot{\beta }}(\eta ^\alpha \sigma _{\alpha \dot{\beta }}^\mu \dot{\overline{\eta }}^{\dot{\beta }}\lambda ^\alpha \sigma _{\alpha \dot{\beta }}^\mu \dot{\overline{\lambda }}^{\dot{\beta }}).$$
Thus the Lagrangian (65) takes in the notation using Dirac spinor $`\psi `$ the following simple form
$$=m(\dot{\omega }_\mu \dot{\omega }^\mu )^{1/2}im(\dot{\overline{\psi }}\psi \overline{\psi }\dot{\psi }),$$
(67)
where
$$\dot{\omega }^\mu =\dot{x}^\mu +i(\dot{\overline{\psi }}\gamma ^\mu \psi \overline{\psi }\gamma ^\mu \dot{\psi }).$$
(68)
We would like to point out that the model with spinorial variables described by Dirac spinor corresponds to the choice of noncompact internal sector, with the metric $`\kappa _{ij}=\mathrm{diag}(1,1)`$. It should be added that the model (67) in different context has been firstly proposed in .
### 4.3 Gupta-Bleuler quantization of the model
The constraints (59)–(61) for $`z=m`$ or equivalently $`\kappa =1`$, written in Dirac notation, are the following
$$Tp^2m^20,$$
(69)
$$D^a\pi ^a+i\overline{\psi }^b(\widehat{p}m)_b{}_{}{}^{a}0,$$
(70)
$$\overline{D}_a\overline{\pi }_ai(\widehat{p}m)_a{}_{}{}^{b}\psi _{b}^{}0.$$
(71)
Here $`\pi ^a`$ and $`\overline{\pi }_a`$ defined as $`\pi ^a=/\dot{\psi }_a`$ and $`\overline{\pi }_a=/\dot{\overline{\psi }}^a`$ are conjugate momenta of $`\psi _a`$ and $`\overline{\psi }^a`$; their Poisson brackets are $`\{\psi _a,\pi ^b\}=\delta _a^b`$ and $`\{\overline{\psi }^a,\overline{\pi }_b\}=\delta _b^a`$. Also we shall use notation $`\widehat{p}\gamma ^\mu p_\mu `$.
From Poisson brackets of the constraints
$$\{\overline{D}_a,D^b\}=2i(\widehat{p}m)_a{}_{}{}^{b},\{D^a,D^b\}=0,\{\overline{D}_a,\overline{D}_b\}=0,$$
(72)
$$\{T,D_a\}=\{T,\overline{D}_a\}=0$$
(73)
we obtain directly that the constraint (69) and the half of the spinorial constraints (70), (71) are first class constraints.
The separation of first and second class spinorial constraints in (70), (71) is achieved by the projectors $`𝒫_\pm \frac{1}{2m}(m\pm \widehat{p})`$ where $`1=(𝒫_++𝒫_{})`$. One can check that on mass shell $`p^2=m^2`$ we obtain $`𝒫_\pm 𝒫_\pm =𝒫_\pm `$, $`𝒫_+𝒫_{}=0`$. From eight real spinorial constraints (70), (71) we construct the following sets of reducible constraints
$$F^a=D^b(\widehat{p}+m)_b{}_{}{}^{a},\overline{F}_a=(\widehat{p}+m)_a{}_{}{}^{b}\overline{D}_{b}^{};$$
(74)
$$G^a=D^b(\widehat{p}m)_b{}_{}{}^{a},\overline{G}_a=(\widehat{p}m)_a{}_{}{}^{b}\overline{D}_{b}^{}.$$
(75)
Due to the relations
$$F^b(\widehat{p}m)_b{}_{}{}^{a}=0,(\widehat{p}m)_a{}_{}{}^{b}\overline{F}_{b}^{}=0;$$
$$G^b(\widehat{p}+m)_b{}_{}{}^{a}=0,(\widehat{p}+m)_a{}_{}{}^{b}\overline{D}_{b}^{}=0$$
on the mass–shell (69) in the set of the constraints $`(F^a,\overline{F}_a)`$ there are only four real independent constraints. Analogously, the constraints $`(G^a,\overline{G}_a)`$ contain as well four real independent constraints. Expressing the constraints (70), (71) in term of the constraints (74), (75) we get
$$D^a=\frac{1}{2m}(F^aG^a),\overline{D}_a=\frac{1}{2m}(\overline{F}_a\overline{G}_a).$$
The constraints (74), (75) satisfy the following Poisson brackets algebra
$$\{\overline{F}_a,F^b\}=2i(\widehat{p}+m)_a{}_{}{}^{b}T,\{F^a,F^b\}=\{\overline{F}_a,\overline{F}_b\}=0,$$
$$\{\overline{F}_a,G^b\}=\{\overline{G}_a,F^b\}=2i(\widehat{p}m)_a{}_{}{}^{b}T,\{\overline{F}_a,\overline{G}_b\}=\{F^a,G^b\}=0,$$
$$\{\overline{G}_a,G^b\}=8im^2(\widehat{p}+m)_a{}_{}{}^{b}2i[2m\delta _a^b+(\widehat{p}+m)_a{}_{}{}^{b}]T,\{G^a,G^b\}=\{\overline{G}_a,\overline{G}_b\}=0.$$
From eight real spinorial constraints present in (70), (71) four independent constraints in $`(F^a,\overline{F}_a)`$ are first class whereas four independent constraints contained in $`(G^a,\overline{G}_a)`$ are second class.
We shall employ the Gupta–Bleuler quantization method by imposing on the wave function all first class constraints ($`T`$, $`F^a`$, $`\overline{F}_a`$) and half of the second class constraints being in involution ($`G^a`$ or $`\overline{G}_b`$). We have two quantizations:
– bosonic chiral quantization, with the wave function satisfying the following wave equations
$$T|\mathrm{\Psi }=0,F^a|\mathrm{\Psi }=0,\overline{F}_a|\mathrm{\Psi }=0,\overline{G}_a|\mathrm{\Psi }=0$$
(76)
– bosonic antichiral quantization with the wave function submitted to the following equations
$$T|\mathrm{\Psi }=0,F^a|\mathrm{\Psi }=0,\overline{F}_a|\mathrm{\Psi }=0,G^a|\mathrm{\Psi }=0.$$
(77)
The reducible constraints $`\overline{F}_a`$ and $`\overline{G}_a`$ are equivalent to primary constraint $`\overline{D}_a`$; similarly the constraints $`F^a`$ and are $`G^a`$ equivalent to $`D^a`$. Therefore one can express the wave equations (76), (77) in other equivalent way
– bosonic chiral quantization:
$$T|\mathrm{\Psi }=0,\overline{D}_a|\mathrm{\Psi }=0,F^a|\mathrm{\Psi }=0$$
(78)
– bosonic antichiral case in which wave function is subjected the following constraints
$$T|\mathrm{\Psi }=0,D^a|\mathrm{\Psi }=0,\overline{F}_a|\mathrm{\Psi }=0.$$
(79)
Let us consider chiral case (78) in more details. Using the realization
$$\pi ^a=i/\psi _a,\overline{\pi }_a=i/\overline{\psi }^a$$
and the momentum-dependent wave function $`\mathrm{\Psi }(p,\psi ,\overline{\psi })`$ one can write down the relations (78) as follows
$$\overline{D}_a\mathrm{\Psi }=i[\frac{}{\overline{\psi }^a}+(\widehat{p}m)_a{}_{}{}^{b}\psi _{b}^{}]\mathrm{\Psi }=0,$$
(80)
$$F^a\mathrm{\Psi }=i\frac{}{\psi _b}(\widehat{p}+m)_b{}_{}{}^{a}\mathrm{\Psi }=0,$$
(81)
$$T\mathrm{\Psi }=(p^2+m^2)\mathrm{\Psi }=0.$$
(82)
The equation (80) has the general solution
$$\mathrm{\Psi }(p,\psi ,\overline{\psi })=e^{\overline{\psi }(\widehat{p}m)\psi }\stackrel{~}{\mathrm{\Psi }}(p,\psi )$$
(83)
where the reduced wave function $`\stackrel{~}{\mathrm{\Psi }}(p,\psi )`$ depends only on $`\psi `$, i. e. we have the expansion
$$\stackrel{~}{\mathrm{\Psi }}(p,\psi )=\underset{n=0}{\overset{\mathrm{}}{}}\psi _{a_1}\mathrm{}\psi _{a_n}\varphi ^{a_1\mathrm{}a_n}(p).$$
(84)
Due to commuting nature of spinor $`\psi _a`$ the component fields $`\varphi ^{a_1\mathrm{}a_n}(p)`$ are totally symmetric
$$\varphi ^{a_1\mathrm{}a_n}(p)=\varphi ^{(a_1\mathrm{}a_n)}(p).$$
(85)
The equations (81) provide the Dirac equations for these fields
$$(\widehat{p}+m)_{a_1}{}_{}{}^{b}\varphi _{}^{a_1a_2\mathrm{}a_n}(p)=0.$$
(86)
We see that the multispinorial fields (85) are Bargman–Wigner fields describing massive particles of spins $`n/2`$. Obviously the Klein–Gordon equation (82) is the consequence of (86).
## 5 Conclusion
The classical $`c`$-number higher spin fields (8586) for any spin are mathematically correct, and provide the relativistic quantum–mechanical description of one–particle states with arbitrary mass and spin (see e. g. ). The concept of bosons and fermions is related with the symmetry properties of multiparticle states, obtained in quantum field theory by quantum fields acting on the vacuum state. The description of higher spin fields presented here (see (84), (85)) does not take into consideration the spin–statistics theorem, however in the framework of first–quantized one-particle classical mechanics we need not to specify the statistics. The transition to the proper spin–statistic relation can be achieved in two way:
i) By introducing classical theory as a suitable limit $`\mathrm{}0`$ of quantized higher spin fields. In such a case the half–integer spin fields will have the Grassmann nature (we recall that fermionic quantum fields are described by infinite–dimensional Clifford algebras which become in the limit $`\mathrm{}0`$ an infinite–dimensional Grassmann algebra).
ii) One can pass from one–particle wave function to the wave function describing multiparticle states by suitable symmetrization procedure (besides bosonic and fermionic multiparticle states one can introduce also parabosonic and parafermionic multiparticle states, with ‘mixed’ symmetry properties).
The wave functions obtained in this paper if used for the description of multiparticle states should be therefore suitably symmetrized: one introduces symmetric products of one–particle wave functions for integer spin fields, and totally antisymmetric products if spin is half–integer. Such a procedure is well-known from the description of multi–particle states in quantum mechanics. If we wish to construct the quantum fields which generate multiparticle states from the vacuum we should multiply the $`c`$-number wave functions by respective bosonic and fermionic creation and annihilation operators. Such a procedure for obtaining fermionic fields with half-integer spin or helicity can be applied to $`N=1`$ massless case (Sect. 3) and $`N=2`$ massive case (Sect. 4), due to the presence of linear field equations.
It should be added that $`c`$-number massive higher spin fields have been obtained also in other papers from different relativistic particle models . We should also add that the realizations of ‘bosonic’ superalgebra was used in for description of physical degrees of freedom of the critical open string with $`N=2`$ conformal symmetry in $`2+2`$ dimensions. Further one can point out that if one introduces fields on twistor spaces (see e. g. ) usually they are also commutative for any spin, or any helicity (in massless case).
#### Acknowledgments
The authors would like to thank E. Ivanov for his interest in this paper and numerous valuable comments. |
warning/0506/hep-lat0506003.html | ar5iv | text | # Topological charge in 1+1 dimensional lattice ϕ⁴ theory
## I Introduction
Topology has an important role in nonperturbative quantum field theory on a lattice. For early work see, for example, earlyref . It is of utmost importance in the dual super-conductivity picture of quark confinement in QCD digiacomo . In field theories with non-trivial topology, different topological sectors are characterized by a conserved topological charge which depends on the boundary behaviour of the theory. This conserved quantity is well-defined in classical continuum field theory.
In 1+1 dimensional classical $`\varphi ^4`$ field theory, there exists a kink solution with a conserved topological charge raja . In the quantum version of the above theory using lattice regularization, past studies on kinks focussed on the calculation of the kink mass ct ; aw . Recently these kinks have also been studied in discrete light cone quantization dlcq .
In this paper we investigate the topological charge in the context of the ordered (broken symmetry) phase of 1+1 dimensional lattice $`\varphi ^4`$ field theory. The short range quantum fluctuations lead to the renormalization of the mass and the coupling constant in the topologically trivial sector of the theory. On the other hand, long range fluctuations lead to a phase transition in the strong coupling region which leads to a breaking of the conservation of topological current and hence the vanishing of the topological charge in the disordered (symmetric) phase. We show that the topological charge in the renormalized theory remains invariant in the ordered phase. To the best of our knowledge, topologocal charge and its invariance in this quantum theory has not been investigated before.
Given that the $`\varphi ^4`$ theory and the Ising model are in the same universality class, one can also define a topolgical charge in the Ising model when considered as a Euclidean quantum field theory in 1+1 dimensions digiacomo2 . However, our aim in this work is to look at the renormalizations of the parameters of the $`\varphi ^4`$ theory so that a sensible definition of topological charge, invariant in the ordered phase of the theory, emerges.
Condensation of the topological objects has been associated with the mechanism of phase transitions in various statistical and quantum field theories digiacomo2 ; sf ; kogut (also see works cited in digiacomo2 ). Our study indicates that condensation of kink-antikinks is a possible mechanism for the order-disorder phase transition in 1+1 dimensional $`\varphi ^4`$ theory.
The plan of this paper is as follows. In Sec. II we present our notation and definition of the topological charge in the bare and renormalized lattice theory. Secs. III, IV and V are devoted to the numerical work with details of various issues that we encounter in the calculation. Finally, Sec. VI contains discussion of our results and conclusion.
## II Topological charge and renormalization
We start from the Lagrangian density in Minkowski space (in usual notation)
$`={\displaystyle \frac{1}{2}}_\mu \varphi ^\mu \varphi {\displaystyle \frac{1}{2}}m^2\varphi ^2{\displaystyle \frac{\lambda }{4!}}\varphi ^4`$ (1)
which leads to the Lagrangian density in Euclidean space
$`_E={\displaystyle \frac{1}{2}}_\mu \varphi _\mu \varphi +{\displaystyle \frac{1}{2}}m^2\varphi ^2+{\displaystyle \frac{\lambda }{4!}}\varphi ^4.`$ (2)
Note that in one space and one time dimensions, the field $`\varphi `$ is dimensionless and the quartic coupling $`\lambda `$ has dimension of mass<sup>2</sup>.
The Euclidean action is
$`S_E={\displaystyle d^2x_E}.`$ (3)
Next we put the system on a lattice of spacing $`a`$ with
$`{\displaystyle d^2x}=a^2{\displaystyle \underset{x}{}}.`$ (4)
Because of the periodicity of the lattice sites in a toroidal lattice the surface terms will cancel among themselves (irrespective of the boundary conditions on fields) enabling us to write
$`(_\mu \varphi )^2=\varphi _\mu ^2\varphi `$ (5)
and on the lattice
$`_\mu ^2\varphi ={\displaystyle \frac{1}{a^2}}\left[\varphi _{x+\mu }+\varphi _{x\mu }2\varphi _x\right].`$ (6)
Introducing dimensionless lattice parameters $`m_0^2`$ and $`\lambda _0`$ by $`m_0^2=m^2a^2`$ and $`\lambda _0=\lambda a^2`$ we arrive at the lattice action in two Euclidean dimensions
$`S={\displaystyle \underset{x}{}}{\displaystyle \underset{\mu }{}}\varphi _x\varphi _{x+\mu }+(2+{\displaystyle \frac{m_0^2}{2}}){\displaystyle \underset{x}{}}\varphi _x^2+{\displaystyle \frac{\lambda _0}{4!}}{\displaystyle \underset{x}{}}\varphi _x^4.`$ (7)
With antiperiodic boundary condition (APBC), we need to identify an order parameter to characterize different phases of the theory since $`\varphi `$ is zero in both symmetric and broken phases. In previous works on kinks in two dimensional lattice field theory ct ; aw , vanishing of the kink mass in the symmetric phase was used to identify the critical coupling. Since the calculation of the kink mass is rather involved, it is preferable to have a simpler choice.
For a $`L^2`$ lattice, we propose $`\overline{\varphi }_{\mathrm{diff}}=\frac{1}{2}\left[\overline{\varphi }_{L1}\overline{\varphi }_0\right]`$ as the order parameter where $`\overline{\varphi }_x\mathrm{kink}\mathrm{g}.\mathrm{s}|\varphi _x|\mathrm{kink}\mathrm{g}.\mathrm{s}`$ with $`|\mathrm{kink}\mathrm{g}.\mathrm{s}`$ denoting the kink ground state. $`\overline{\varphi }_x`$ is computed by taking the average of $`\varphi _x`$ over configurations with APBC in the spatial direction (which effectively performs importance sampling around the kink configurations in a Monte-Carlo simulation).
In the classical theory, the topological charge is given by
$`Q_{\mathrm{classical}}=\sqrt{{\displaystyle \frac{\lambda }{6m^2}}}\varphi _{\mathrm{diff}}=\sqrt{{\displaystyle \frac{\lambda }{3m_B^2}}}\varphi _{\mathrm{diff}}`$ (8)
where $`\varphi _{\mathrm{diff}}=\frac{1}{2}\left[\varphi (\mathrm{})\varphi (\mathrm{})\right]`$ and we have used the fact that the mass of the elementary boson is $`m_B=\sqrt{2}m`$ in the broken phase. In the classical theory, the relation between $`\varphi _{\mathrm{diff}},m_B`$ and $`\lambda `$ is given by
$`\varphi _{\mathrm{diff}}=\sqrt{{\displaystyle \frac{3m_B^2}{\lambda }}}`$ (9)
which guarantees that the topological charge is +1 (kink sector) and -1 (antikink sector) in the broken symmetric phase and zero in the symmetric phase.
In the lattice theory with APBC in the $`\mathrm{`}\mathrm{`}`$spatial” direction, the bare topological charge is defined by
$`Q_0=\sqrt{{\displaystyle \frac{\lambda _0}{3m_{B_0}^2}}}\overline{\varphi }_{\mathrm{diff}}.`$ (10)
Due to quantum fluctuations, the classical relation analogous to Eq. (9) between $`\overline{\varphi }_{\mathrm{diff}},m_{B_0}`$ and $`\lambda _0`$ is not obeyed and as a consequence, $`Q_0`$ does not remain +1 or -1 in the ordered phase.
We propose to define the topological charge in the renormalized theory as
$`Q_R=\sqrt{{\displaystyle \frac{\lambda _R}{3m_R^2}}}\overline{\varphi }_{R_{\mathrm{diff}}}`$ (11)
where $`\varphi _R=\frac{1}{\sqrt{Z}}\varphi `$, $`Z`$ being the field renormalization constant. The conventional definition of the renormalized coupling $`\lambda _R`$ in $`\varphi ^4`$ theory is in terms of the renormalized four-point vertex function. The calculation of the four-point vertex function on the lattice is computationally demanding. Fortunately, in the broken phase, we can choose a definition munster of $`\lambda _R`$ in terms of the renormalized mass $`m_R`$ and the renormalized vacuum expectation value $`\varphi _R`$, determined with periodic boundary condition (PBC):
$`\lambda _R=3{\displaystyle \frac{m_R^2}{\varphi _R^2}}.`$ (12)
This definition of $`\lambda _R`$ involves only the computation of one-point function for $`\varphi `$ and the two point function for $`m_R`$ and $`Z`$. The details of these computations are provided in Ref. longpaper . Using this definition of $`\lambda _R`$ in Eq. (11), we get
$`Q_R={\displaystyle \frac{\overline{\varphi }_{R_{\mathrm{diff}}}}{\varphi _R}}={\displaystyle \frac{\overline{\varphi }_{\mathrm{diff}}}{\varphi }}.`$ (13)
The second equality of Eq. (13) assumes $`Z`$ to be the same for $`\varphi _{\mathrm{diff}}`$ and $`\varphi `$ and therefore involves two unrenormalized quantities, $`\overline{\varphi }_{\mathrm{diff}}`$ calculated with APBC and $`\varphi `$ calculated with PBC. If $`\overline{\varphi }_{\mathrm{diff}}`$ and $`\varphi `$ are of the same magnitude, the renormalized topological charge is $`\pm 1`$ in the ordered phase except for the statistical and systematic errors that occur in the numerical computations of the quantities in the numerator and the denominator of Eq. (13).
## III Calculation with PBC: $`\varphi `$
First we discuss the calculation of $`\varphi `$ with PBC. As is well-known, spontaneous symmetry breaking cannot occur at finite volume due to the quantum mechanical tunneling phenomena. This is due to the fact that at finite volume we have finite degrees of freedom. Tunneling probability is, however, inversely proportional to the exponential of the volume and hence is exponentially suppressed at large volume munster ; montvay . Thus in practice, the issue is how large a volume we need for the suppression to occur.
For $`\varphi ^4`$ theory, tunneling is present for $`10^2`$ lattice but is practically absent for $`24^2`$ lattice. This is demonstrated in Fig. 2 where $`\varphi `$ averaged over all sites for a given configuration is plotted against the configuration number. These calculations were done using single hit Metropolis algorithm.
We find $`\varphi `$ to be sensitive to the various choices of the initial configuration (see Fig. 2) in the simulation using the Metropolis algorithm. For cold starts \[$`\varphi `$ = +1 (or -1) for all sites\] the selection of the vacuum is controlled by the choice of the sign of $`\varphi `$ irrespective of the couplings. On the other hand for hot starts ($`\varphi =\pm 1`$ sprinkled randomly over the entire lattice), the vaccum is picked randomly from the two degenerate vacua.
A more serious problem with the Metropolis algorithm is critical slowing down which makes computation near the critical point prohibitively expensive. In the Monte-Carlo simulation, to reduce critical slowing down, we have thus resorted to the standard Metropolis update combined with a cluster algorithm wolff to update the embedded Ising variables in $`\varphi ^4`$ theory, following the method of Brower and Tamayo bt . We use Wolff’s single cluster variant of the algorithm. This method has been used previously in determining the phase diagram of two dimensional lattice $`\varphi ^4`$ theory lw .
In the cluster algorithm however, since the sign of the field of all the members of the cluster are flipped in every updation cycle, the algorithm actually enforces tunneling and the configuration average of $`\varphi `$ is always zero. Thus to get the appropriate nonzero value for the condensate we measure $`|\varphi |`$ where $`\varphi =\frac{1}{\mathrm{Volume}}\underset{\mathrm{sites}}{}\varphi (x)`$. To understand the mod let us consider a local order parameter $`\varphi (x)`$. Since the configurations will be selected at random dominantly from the neighborhood of either vacua in the broken phase, $`\varphi (x)`$ will vanish when averaged over configurations thus wiping out the signature of a broken phase. If one uses $`|\varphi (x)|`$ as the order parameter then in the broken phase it correctly projects itself onto one of the vacua yielding the appropriate non-zero value. The use of this mod, unfortunately, destroys the signal in the symmetric phase completely by wiping out the significant fluctuations in sign. However if we choose to use $`|\frac{1}{\mathrm{Volume}}\underset{\mathrm{sites}}{}\varphi (x)|`$, it correctly captures the broken phase as well as the symmetric phase. While the sign fluctuation over configurations are still masked, the fluctuations over sites survive producing $`|\varphi |=0`$ correctly in the symmetric phase.
The volume dependence of $`|\varphi |`$ is presented in Fig. 4 where we also compare with the classical value of the condensate. As can be seen, for small volumes, the signal for phase transition is very weak to detect, but $`50^2`$ lattice is big enough to observe the transition.
In Figs. 2 and 4 and all the figures to follow the standard error bars, if not visible, are smaller than the symbols unless otherwise stated.
The phase diagram we obtained for a $`512^2`$ lattice using PBC is presented in Fig. 4. This agrees with the phase diagram obtained in ct ; aw ; lw . In lw the authors extrapolate their results to infinite volume. We observe that our $`512^2`$ lattice results are as good as the infinite volume result in Ref. lw .
Most of our calculations are performed at $`m_0^2=0.5`$. The corresponding critical coupling $`\lambda _0^c`$ is around 1.95.
## IV Calculation with APBC: Kink configurations and $`\overline{\varphi }_{\mathrm{diff}}`$
Cluster algorithms are known to fail with antiperiodic boundary conditions in $`\varphi ^4`$ theory hasen . For the calculation of $`\overline{\varphi }_{\mathrm{diff}}`$, because of APBC, we cannot use the cluster algorithm and we resort to the standard Metropolis algorithm. In addition to the problems associated with critical slowing down, we also face problems associated with computing the profile of an extended object, the topological kink in this case and also in fixing its location for the measurement of topological charge.
First we note that the location of the kink cannot be controlled with a cold start as we show in Fig. 6. It is perhaps not easy at first sight, to recognize that most of the configurations seen in Fig. 6 are actually kink configurations. To bring it out we refer to Fig. 8 in which we demonstrate schematically, how a kink centered near the boundary of a toroidal lattice with APBC on the fields should look like. Note that the $`(L+i)`$ th site is identified with the $`i`$ th site on a toroidal lattice but the sign of the field is flipped on account of APBC on the fields giving it a plateau-like appearance. Most of the configurations in Fig. 6 resemble this plateau.
All such profiles in Fig. 6 are therefore basically kinks or antikinks located near the edge. We think that the formation of kinks near the edge are favored by the algorithm because it is clearly much easier to generate the plateau-like configurations from a cold start ($`\varphi =+1`$) compared to a kink which involves a flipping of the sign of fields over larger region. Near the critical region of course, we expect the formation of kink like configurations to be much easier and it is consistent with our observation ($`\lambda _0=1.8`$ in Fig. 6). Let us mention at this point that because of translational invariance, the position of a kink is of no significance as long as it does not tend to the spatial infinities.
To make our results more transparent we intended to work with kink configurations that are centered near the middle of the lattice. The definition of the topological charge that we use (Eq. (11)) also presupposes that our kink is actually located near the center.
To obtain such kink configurations we generated configurations by using the kink start ($`\varphi _x=1`$ for $`0x<\frac{L}{2}`$ and $`\varphi _x=+1`$ for $`\frac{L}{2}x<L`$) which nicely generates the kink configurations near the center (Fig. 6).
In passing we remark that the data in both, Fig. 6 and Fig. 6 are seen to deteriorate with smaller volumes owing to finite size effects which are as usual more serious near the critical point (large coupling in this case). We have also observed that there is a small but noticeable increase in the size of the kinks (spatial extent over which $`\varphi _x`$ changes) with an increase of $`\lambda _0`$ which affects the sharpness of the kinks in a small volume near the critical region.
Reasonably good kink configurations for the measurement of topological charge (with stable flat regions) are obtained with the kink start with $`128^2`$ and $`256^2`$ lattices.
We have observed kink motion in the vicinity of the critical region. This motion of the kink is clearly visible in Fig. 8 where we present $`\overline{\varphi }_x`$ for a $`512^2`$ lattice for three different time slices. Probably because the kink mass is small enough at the couplings, we see this motion. This is consistent with the vanishing of the kink mass at the critical point. Incidentally, the equality in the magnitude of $`\overline{\varphi }_{\mathrm{diff}}`$ among different time slices demonstrates the expected conservation of topological charge. One would have ideally liked to increase the statistics by averaging over the time slices at a given $`x`$, but the movement of kink makes this not viable in general. However, taking into account the fact that the location of the kink is of no consequence in the measurement of topological charge, we have actually forced the formation of kinks at the middle (by fixing the value of the field at the mid-point in every time slice, to zero) and then taken the average over time slices.
## V Kink condensation and topological charge
Let us now come to the discussion of the topological charge. For the use of Eq. (13) for the topological charge, it is necessary that we get the same phase diagram and critical behaviour for the numerator and the denominator of this expression. The simulation results with PBC and APBC shown in Fig. 10 for $`48^2`$ lattice indeed confirm this within our numerical accuracy. Here we would also like to remark that Fig. 10 shows that the disappearance of the kink configuration (characterized by $`|\overline{\varphi }_{\mathrm{diff}}|=0`$) and the onset of the phase transition from broken to symmetric phase (characterized by $`|\varphi |=0`$) coincide. It therefore suggests ‘kink-condensation’ as a possible mechanism for the order-disorder phase transition in 1+1 dimensional $`\varphi ^4`$ theory. Such a connection is known to exist in 2-dimensional Ising model sf ; kogut ; digiacomo2 which is believed to be in the same universality class and is not totally unexpected in $`\varphi ^4`$ theory since it has an embedded Ising variable. Condensation of topological excitation has generally been associated with the mechanism of phase transition in many statistical and quantum field theories. Kink condensation is clearly visible in Fig. 10 where we present $`\overline{\varphi }_x`$ for three couplings very close to the critical region. At $`\lambda _0=1.95`$, kink configuration is just barely visible, i.e., the boundary value of $`\overline{\varphi }_x`$ is still nonzero and $`\overline{\varphi }_x`$ passes through zero only once. As the coupling increases, the boundary value of $`\overline{\varphi }_x`$ reduces to zero signalling disappearance of the single kink. However, $`\overline{\varphi }_x`$ passes through zero many times showing closely packed kink-antikink configurations of varying amplitudes.
Finally we present the result for the topological charge in Figs. 12 and 12. The behavior of the topological charge for a range of $`\lambda _0`$ is depicted in Fig. 12 and Fig. 12 for $`48^2`$ and $`512^2`$ lattices respectively. The figures clearly demonstrate the need for renormalization: It restricts $`Q_R`$ to +1 for the whole range of $`\lambda _0`$ investigated in conformity with our expectations. Near the critical region the fluctuation of $`Q_R`$ around unity is due to different systematic and statistical errors associated with different numerical algorithms used for the calculation of the numerator and the denominator of Eq. (13) especially because the numerator is evaluated with Metropolis algorithm which is known to suffer from critical slowing down. We also see by comparing Figs. 12 and 12 that there is noticeable improvement of data with volume near the critical region.
## VI Summary and conclusions
In this work we have investigated the topological charge in 1+1 dimensional lattice $`\varphi ^4`$ field theory. We have shown that with APBC in the spatial direction, lowest energy configuration is a kink or an antikink. In order to characterize the different phases of the theory with APBC we have proposed a simple order parameter $`\overline{\varphi }_{\mathrm{diff}}`$ and we have demonstrated its effectiveness.
In the process of computing the topological charge we have observed that as the system moves on from ordered to disordered phase, the single kink (or antikink) gives way to the occurance of a multitude of kink-antikink pairs suggesting kink condensation as a possible mechanism for this phase transition.
A major issue in extracting the topological charge is the renormalization. By a particular choice of the renormalized coupling (in the topologically trivial sector), we are able to express the topological charge in the renormalized theory as the ratio of renormalized order parameters in the lattice theories with APBC and PBC. Provided the field renormalization constants are the same in the two cases, the expression for the topological charge in the renormalized theory becomes the ratio of unrenormalized order parameters. Making use of this, we have computed the topological charge in the renormalized theory and demonstrated that it indeed maintains the value +1 in the broken phase.
###### Acknowledgements.
Numerical calculations presented in this work are carried out on a Power4-based IBM cluster and a Cray XD1. The High Performance Computing Facility is supported by the 10<sup>th</sup> Five Year Plan Projects of the Theory Division, Saha Institute of Nuclear Physics, under the DAE, Govt. of India. |
warning/0506/cond-mat0506414.html | ar5iv | text | # Design and control of spin gates in two quantum dots arrays
## Abstract
We study the spin-spin interaction between quantum dots coupled through a two dimensional electron gas with spin-orbit interaction. We show that the interplay between transverse electron focusing and spin-orbit coupling allows to dynamically change the symmetry of the effective spin-spin Hamiltonian. That is, the interaction can be changed from Ising-like to Heisenberg-like and vice versa. The sign and magnitude of the coupling constant can also be tuned.
Coherent control and measurement of quantum spins are at the heart of new technologies with great potential value for information processing Awschalom et al. (2002); Nielsen and Chuang (2000). This has lead to a great activity in the field of quantum spin control in solid state devices Loss and DiVincenzo (1998); Burkard et al. (1999); Hu and Das Sarma (2000); Divicenzo et al. (2000); Potok et al. (2002); Elzerman et al. (2003); Kato et al. (2004, 2005). Since the seminal work by Loss and DiVincenzo Loss and DiVincenzo (1998), the exchange gate is the central tool underlying most of the proposals for spin manipulation in solid state devices based on quantum dots (QDs) Burkard et al. (1999); Hu and Das Sarma (2000); Divicenzo et al. (2000). The exchange gate is founded on the Heisenberg interaction between localized spins and so far nearly all implementations for such control are based on an ‘on/off’ setup—the interaction is either active or inactive. Furthermore, when the exchange gate is controlled by electrical gates, the control implies to ‘open’ or ‘close’ the QDs, changing their coupling to the environment, their shape and thus their detailed internal electronic structure. The question then is : Is it possible to engineer a predefined spin-spin interaction between QDs and then change its magnitude, sign and symmetry with a negligible impact on the internal structure of the dots?
In a recent work, we analyzed a way to tune the amplitude and sign of the spin coupling Usaj et al. (2005). Here we go a step forward and show how to design a Heisenberg or an Ising-like interaction of the desired magnitude and sign of the coupling constant and then dynamically change one into the other by controlling a small magnetic field—the control mechanism relies on the interplay between transverse electron focusing and spin-orbit coupling van Houten et al. (1989); Beenakker and van Houten (1991); Usaj and Balseiro (2004); Rokhinson et al. (2004). This opens up the possibility to manipulate spin-spin Hamiltonians in solid state devices as it is done today with NMR techniques in molecules Ernst et al. (1990).
The setup consists of two QDs at the edges of two electron gases as schematically shown in Fig. 1a, with an interdot distance $`d`$ of the order of $`1\mu `$m. Present semiconducting heterostructure technology allows tailoring this structure in two dimensional electron gases (2DEG). In the Coulomb blockade regime, the QDs can be gated to have an odd number of electrons so that they behave as magnetic objects. In what follows we describe them as localized $`\frac{1}{2}`$ spins. The virtual tunneling of electrons between the dots and the 2DEG leads to a Kondo coupling between the localized spins $`\stackrel{}{𝑺}_i`$ and the 2DEG spins described by the following Hamiltonian Schrieffer and Wolf (1966):
$$\widehat{H}_K=\underset{i,\eta ,\eta ^{}}{}J_i\stackrel{}{𝑺}_i\psi _{\eta \sigma }^{}(R_i)\frac{\stackrel{}{𝝈}_{\sigma \sigma ^{}}}{2}\psi _{\eta ^{}\sigma ^{}}(R_i)$$
(1)
where $`i=1,2`$ indicates the left and right QD respectively, $`\psi _{\eta \sigma }^{}(R_i)`$ creates an electron with spin $`\sigma `$ in a Wannier-like orbital centered around the coordinate $`R_i`$ of the $`i`$-th QD at the upper ($`\eta =1`$) or lower ($`\eta =2`$) plane. The spacial extension of the Wannier orbital depends on the opening of the QDs. This coupling leads to a RKKY-like interaction between the QDs spins that takes the general form Imamura et al. (2004):
$$\widehat{H}_J=\frac{J_1J_2}{4\pi }\mathrm{Im}𝑑\omega f(\omega )\mathrm{Tr}\left(\stackrel{}{𝑺}_1\stackrel{}{𝝈}𝑮(1,2)\stackrel{}{𝑺}_2\stackrel{}{𝝈}𝑮(2,1)\right)$$
(2)
where $`f(\omega )`$ is the Fermi function and the $`2\times 2`$ matrix $`𝑮(i,j)`$ is the Fourier transform of the retarded electron propagator whose elements are $`G_{\sigma \sigma ^{}}(i,j,tt^{})=\mathrm{i}\theta (tt^{})\times `$ $`_{\eta ,\eta ^{}}\{\psi _{\eta \sigma }(R_i,t),\psi _{\eta ^{}\sigma ^{}}^{}(R_j,t^{})\}`$. When the electron’s spin is conserved along the electron propagation between QDs, $`𝑮(i,j)`$ is diagonal in the spin index and the spin-spin Hamiltonian (2) reduces to the Heisenberg one $`\widehat{H}_J=J\stackrel{}{𝑺}_1\stackrel{}{𝑺}_2`$ with the usual RKKY-like exchange $`J=J_1J_2/2\pi \mathrm{Im}𝑑\omega f(\omega )G_{}(1,2)G_{}(2,1)`$.
The presence of a small magnetic field $`B_z`$ perpendicular to the 2DEG creates edge states that dominate the electron scattering from objects placed at the 2DEG edges. The interaction between QDs is then mediated by these edge states and the propagators are mainly due to the semiclassical orbits shown in Figs. 1b and 1c. Due to the chiral nature of these orbits, the intra-plane scattering, described by the terms in Eq. (1) with $`\eta =\eta ^{}`$, give forward scattering while the inter-plane terms (with $`\eta \eta ^{}`$) describe the backward scattering. Only the inter-plane backward scattering processes contribute to the effective interaction. In other words, each propagators in Eq. (2) is due to contributions from only one plane. As the external field increases the cyclotron radii of these orbits decrease: $`r_c=\mathrm{}kc/eB_z`$ with $`k`$ the electron wavevector. The focusing fields are those for which the interdot distance $`d`$ is commensurate with the cyclotron radius $`r_c`$ of electrons at the Fermi energy ($`E_\text{F}`$), that is $`d=2nr_c=2n`$ $`\mathrm{}k_\text{F}c/eB_z`$ with $`n`$ an integer number. At the focusing fields, the electrons at the Fermi level scattered by one QD are focused onto the other leading to an amplification of the exchange integral $`J`$ . The numerical result is shown in Fig. 1d where, for the sake of comparison with the conventional RKKY interaction, the exchange integral $`J`$ is plotted as a function of the interdot distance for a fixed magnetic field. These results were obtained using a finite differences technique Usaj and Balseiro (2004) for a system with an effective electronic mass $`m^{}=0.067m_e`$ and $`E_\text{F}=5`$meV, corresponding to an electron density of $`1.5\times 10^{11}/cm^2`$. With these parameters, the focusing amplification of the exchange integral is clearly observed. In the semiclassical picture, the first focusing condition ($`n=1`$) corresponds to a direct propagation of the electrons from one QD to the other; in the second one ($`n=2`$) the electron bounces once at the 2DEG edge. For interdot distances of the order of $`1\mu `$m, the magnetic fields for the first focusing conditions ($`n=1`$ or $`2`$) is small and neglecting the Zeeman spin splitting due to the external field is a good approximation. It is worth mentioning that, in similar geometries, the electron focusing due to small magnetic fields is clearly observed in transport experiments van Houten et al. (1989); Potok et al. (2002); Folk et al. (2003).
In systems with strong spin-orbit (SO) coupling, new effects arise. We consider a Rashba SO interaction in the 2DEG Rashba (1960); Bychkov and Rashba (1984). This interaction is due to the inversion asymmetry of the confining potential and it is described by the Hamiltonian $`H_{SO}=\alpha /\mathrm{}(p_y\sigma _xp_x\sigma _y)`$ where $`p_\gamma `$ are the components of the canonical momentum of the 2DEG electrons and $`\sigma _\gamma `$ the spin operators. The SO coupling acts as a strong in-plane magnetic field proportional to the momentum. This breaks the spin degeneracy leading to two different conduction bands Rashba (1960).
In the presence of a small magnetic field perpendicular to the gas plane, each band leads to a different cyclotron radius. These two radii manifest as two distinct focusing fields for the first ($`n=1`$) focusing condition Usaj and Balseiro (2004). This splitting has been observed by Rokhinson et al. Rokhinson et al. (2004) in a p-doped GaAs/AlGaAs heterostructure. The spin texture of the orbits is such that, for small fields (large cyclotron radii), the electron’s spin adiabatically rotates along the semiclassical orbit, being perpendicular to the momentum, as schematically shown in Fig. 2. In order to describe the magnetic scattering of electrons in this case, it is convenient to quantize the spin along the $`x`$-axis. Then, around the first focusing condition the propagators $`G_{+,}(i,j)`$ and $`G_{,+}(i,j)`$ dominate the interdot coupling, here the spin index $`\pm `$ indicate the two spin projection. The interdot interaction is then approximately given by an Ising term $`\widehat{H}_I=J_{xx}S_{1x}S_{2x}`$ with coupling constant given by
$$J_{xx}=\frac{J_1J_2}{4\pi }\mathrm{Im}𝑑\omega f(\omega )G_{+,}(i,j)G_{,+}(j,i)$$
(3)
where $`i=1`$ and $`j=2`$ or $`i=2`$ and $`j=1`$ depending on which cyclotron radius contributes to the focusing. This result can be visualized in terms of the semiclassical trajectories shown in Figs. 2a and 2b: for a SO coupling strong enough to split the focusing condition, the inter-plane spin-flip backscattering mixes the two cyclotron radii living the electron out of the focusing condition. Thus, these spin-flip processes can not contribute to the coupling. The interdot interaction is then due to non-spin flip processes of electrons that are back-scattered. This defines the symmetry axis of the resulting Ising interaction.
At the second focusing condition, the system operates in a different way (see Fig. 2c). There are two important effects to consider: i) the orbits with different cyclotron radii are mixed at the bouncing point due to spin conservation, and ii) along the trajectories from one QD to the other the electron’s spin completes a $`2\pi `$ rotation. As a consequence, the two orbits contribute to the exchange integral and $`G_{+,+}(i,j)`$ and $`G_,(i,j)`$ dominate the spin-spin coupling. In this way the rotational symmetric Heisenberg coupling is recovered.
For arbitrary external field, Hamiltonian (2) can be written as fully anisotropic Heisenberg model plus a Dzyaloshinski-Moriya term
$$\widehat{H}_J=\underset{\gamma }{}J_{\gamma \gamma }S_{1\gamma }S_{2\gamma }+\stackrel{}{\beta }\left(\stackrel{}{𝑺}_1\times \stackrel{}{𝑺}_2\right)$$
(4)
where $`\stackrel{}{\beta }=(0,\beta _0,0)`$. Hamiltonian (4) is a particular case of a more general Hamiltonian including SO effects Dzyaloshinski (1958); Moriya (1960); Bonesteel et al. (2001). In our case, due to the symmetry of our geometry, there are only four independent parameters: $`J_{\gamma \gamma }`$ with $`\gamma =x,y,z`$ and $`\beta _0`$. The numerical results for these coupling constants are shown in Fig. 2d. As argued above, around the first focusing condition the system behaves as an Ising like model: the dominant coupling $`J_{xx}`$ shows a large amplification when the interdot distance matches each one the two cyclotron orbits. The relative amplitude and sign of $`J_{xx}`$ in these peaks depends on both the external field and the Fermi energy. At the second focusing condition the system behaves as an isotropic Heisenberg model ($`J_{xx}=J_{yy}=J_{zz}`$) with a small anisotropic correction ($`|\beta _0/J_{xx}|1`$). Figures 3a and 3b show the dominant couplings, $`J_{xx}`$ and $`J_{zz}`$ respectively, as a function of $`B_z`$ and $`d`$. The magnetic field not only can turn on and off each coupling but a fine tune around the focusing fields can change their sign too (see Figs. 3c and 3d) Usaj et al. (2005).
There is a variety of systems that are potentially appropriate to observe these effects. While in n-doped GaAs/AlGaAs heterostructures the spin orbit is small, systems like p-doped GaAs/AlGaAs or InGaAs heterostructures present a large SO coupling. The nature of the SO effect depends on the system. Effects like the ones discussed in this work are also present in systems with Dresselhaus SO coupling. Furthermore, the external control of the relative magnitude of both contributions to the SO coupling Miller et al. (2003), could allow the control the quantization axis of the Ising-like interaction.
In summary, we have shown that the interplay between transverse electron focusing and spin orbit interaction gives a unique opportunity to control and tune the spin-spin interaction between QDs without inducing big changes in their internal structure. When the SO coupling is large, it leads to a spin-dependent focusing condition (for $`n=1`$), resulting in a highly anisotropic Ising-like interaction. However, by doubling the external field a Heisenberg gate with a small correction of the Dzyaloshinski-Moriya type is recovered. In the context of quantum computing, there are strategies to eliminate or control these small corrections to the Heisenberg gate Bonesteel et al. (2001); Stepanenko et al. (2003); Stepanenko and Bonesteel (2004). The proposed setup can be extended to three or more QDs in a linear array. An array with different interdot distances and with extra gates used to blockade the focusing may be used to independently control the interdot interactions.
This work was partially supported by ANPCyT Grants No 13829 and 13476 and Fundación Antorchas, Grant 14169/21. GU acknowledge support from CONICET. |
warning/0506/math0506275.html | ar5iv | text | # The Geometric Theory of the Fundamental Germ
## 1. Introduction
This paper represents a continuation of our quest to extend $``$-coefficient algebraic topology to laminations through the generalization of $`\pi _1`$ called the fundamental germ. In this paper, we extend this construction to any lamination admitting a smooth structure.
Let us recall briefly the intuition behind the fundamental germ. Consider a suspension
$$_\rho =\left(\stackrel{~}{B}\times 𝖳\right)/\pi _1B$$
of a representation $`\rho :\pi _1B\mathrm{𝖧𝗈𝗆𝖾𝗈}(𝖳)`$, where $`B`$ is a manifold. Then $`\pi _1B`$ acts on $`_\rho `$ as fiber preserving homeomorphisms. Let $`T𝖳`$ be a fiber transversal and let $`x_0,x,T`$. A $`\pi _1B`$-diophantine approximation of $`xT`$ based at $`x_0`$ is a sequence $`\{g_\alpha \}\pi _1B`$ with $`g_\alpha x_0x`$. The fundamental germ $`[[\pi ]]_1(,x_0,x)`$ is then the groupoid of tail equivalence classes of sequences of the form $`\{g_\alpha h_\alpha ^1\}`$ where $`\{g_\alpha \}`$, $`\{h_\alpha \}`$ are diophantine approximations of $`x`$ along $`x_0`$ . This construction is more generally available for any lamination occurring as a quotient of a suspension, a double-coset of a Lie group or a locally-free action of a Lie group on a space, laminations which we refer to collectively as algebraic. Intuitively, if $`L`$ is the leaf containing $`x_0`$, the elements of $`[[\pi ]]_1(,x_0,x)`$ can be thought of as sequences of paths in $`L`$ whose endpoints converge transversally to $`x`$. Such a sequence can be thought of as an ideal loop based at $`x`$ that records an “asymptotic identification” within the leaf $`L`$.
For a linear foliation $`_r`$ of a torus by lines of slope $`r`$, the diophantine analogy is literal and $`[[\pi ]]_1(,x_0,x)`$ is the group of classical diophantine approximations of $`r`$. A manifold $`B`$ is a supension of the trivial representation i.e. a lamination with a single leaf and fiber transversals that are points, which forces $`x_0=x`$. Then all sequences in $`\pi _1B`$ converge, and we find that $`[[\pi ]]_1(B,x)={}_{}{}^{}\pi _{1}^{}(B,x)=`$ the nonstandard fundamental group of $`B`$.
We now turn to the contents of this article. The algebraic definition of the fundamental germ just described, while amenable to calculation, has the following serious drawbacks:
1. It is available only for the select family of algebraic laminations.
2. It is an invariant only with respect to the special class of trained lamination homeomorphisms (c.f. ).
Addressing these flaws is the central theme of the present study. In the summary that follows, we shall assume for simplicity that all leaves are simply connected.
We begin with item (1). Let $``$ be an arbitrary lamination admitting a smooth structure, let $`x_0,x`$ be as above and denote by $`L`$ the leaf containing $`x_0`$. Equip $``$ with a leaf-wise riemannian metric that has continuous transverse variation. In this paper, we shall refer to such a lamination as riemannian. The new idea here is to use the leaf-wise geometry to represent – as sequences of isometries – the diophantine approximations which would make up $`[[\pi ]]_1`$. If $`L`$ has constant curvature geometry, this prescription may be followed word-for-word. Fixing a transversal $`T`$ containing $`x_0,x`$ and a continuous section of orthonormal frames $`𝖿=\{𝖿_y\}`$, $`yT`$, we define a diophantine approximation of $`x`$ to be a sequence $`\{A_\alpha \}`$ of isometries of $`L`$ for which $`(A_\alpha )_{}𝖿_{x_0}`$ belongs to $`𝖿`$ and converges transversally to $`𝖿_x`$. The fundamental germ $`[[\pi ]]_1(,x_0,x,𝖿)`$ is then defined to be the set of tails of sequences of the form $`\left\{A_\alpha B_\alpha ^1\right\}`$ where $`\{A_\alpha \}`$, $`\{B_\alpha \}`$ are diophantine approximations of $`x`$.
In the case of non constant curvature leaf-wise geometry, it is necessary to work within the category of virtual geometry in order to make sense of the notion of diophantine approximation. There, a riemannian manifold $`M`$ is replaced by a union of riemannian manifolds, its virtual extension $`{}_{}{}^{}M`$, which consists of all sequences in $`M`$ up to the relation of being asymptotic. A virtual isometry between riemannian manifolds $`M`$ and $`N`$ consists of a pair of isometric inclusions $`{}_{}{}^{}M{}_{}{}^{}N`$. All dense leaves of a riemannian lamination have virtually isometric universal covers, and moreover, a dense leaf having no ordinary isometries will admit many virtual isometries.
This leads to the following definition of a diophantine approximation: let $`xT`$, $`𝖿`$ a frame field on $`T`$ and let $`L`$ be any leaf accumulating on $`x`$. Then a sequence $`𝖿_{x_\alpha }𝖿_x`$, $`\{x_\alpha \}L`$, determines an isometry $`{}_{}{}^{}f:L_xU{}_{}{}^{}L`$, where $`L_x`$ is the leaf containing $`x`$ and $`U`$ is a component of $`{}_{}{}^{}L`$. The fundamental germ $`[[\pi ]]_1(,L,x,𝖿)`$ is defined to be the set of (maximal extensions of) maps of the form
$${}_{}{}^{}f{}_{}{}^{}g_{}^{1}.$$
In this way, we now have a definition of the fundamental germ valid for any lamination admitting a smooth structure along the leaves.
In order to address drawback (2), we will need the germ universal cover
$$[[\stackrel{~}{}]]{}_{}{}^{}L,$$
defined to be the set of asymptotic classes of sequences in $`L`$ that converge to points of $``$. The germ universal cover plays the role of a unit space for a groupoid structure on $`[[\pi ]]_1(,L,x,𝖿)`$. It is a lamination whose leaves are nowhere dense, and when $`L`$ is dense, the canonical map $`[[\stackrel{~}{}]]`$ is onto. We may therefore think of $`[[\stackrel{~}{}]]`$ as obtained from $``$ by “unwrapping” all transversal topology implemented by $`L`$.
Assume now that $`L`$ is dense. The mother germ $`[[\pi ]]_1`$ is defined to be the groupoid of all partially defined maps of $`[[\stackrel{~}{}]]`$ that are homeomorphisms on domains which are sublaminations of $`[[\stackrel{~}{}]]`$ and preserve the projection $`[[\stackrel{~}{}]]`$. We have in particular that
$$[[\pi ]]_1\backslash [[\stackrel{~}{}]].$$
$`[[\pi ]]_1`$ is the receptacle of all the $`[[\pi ]]_1(,L^{},x,𝖿)`$ for $`L^{}`$ dense, in that it contains subgroupoids isomorphic to each. The mother germ is functorial with respect to topological lamination covering maps, and is therefore, in spite of its riemannian construction, a topological invariant. This takes care of item (2) above.
The remainder of the paper is devoted to examples and an application. Many examples were discussed in , and so for this reason we limit ourselves to laminations which are not algebraic and hence which do not have a fundamental germ in the sense described there.
The first example we consider is that which we call here the antenna lamination, a surface lamination discovered by Kenyon and Ghys which has the distinction of having leaves of both parabolic and hyperbolic type. With respect to a hyperbolic leaf, the fundamental germ is calculated as a set to be $`{}_{}{}^{}F_{2}^{}\times ({}_{}{}^{}_{\widehat{2}}^{}{}_{}{}^{}_{\widehat{2}}^{})`$ where $`{}_{}{}^{}F_{2}^{}`$ is the nonstandard free group on two generators, and $`{}_{}{}^{}_{\widehat{2}}^{}`$ is the subgroup of $`{}_{}{}^{}`$ isomorphic to the fundamental germ of the dyadic solenoid. Although a product of groups, this germ is not a group with respect to its defined multiplication. It is the first example we have encountered of a fundamental germ that is not a group.
The second example is that of the Anosov foliation of the unit tangent bundle to the modular surface. Although this is just the suspension of the action of $`PSL(2,)`$ on the boundary of the hyperbolic plane, the definition of the fundamental germ found in is unavailable since it does not work for actions with fixed points. We calculate the fundamental germ here as a set to be $`PSL(2,{}_{}{}^{})`$, but as in the case of the antenna lamination, it is also not a group with respect to its defined multiplication.
The final result of this paper concerns the use of the fundamental germ to calculate the mapping class group of the algebraic universal cover $`\widehat{\mathrm{\Sigma }}`$ of a closed surface $`\mathrm{\Sigma }`$ of hyperbolic type. $`\widehat{\mathrm{\Sigma }}`$ is by definition the inverse limit of finite covers of $`\mathrm{\Sigma }`$, a compact solenoid with dense disk leaves. If $`L\widehat{\mathrm{\Sigma }}`$ is a fixed leaf, the leafed mapping class group $`\mathrm{𝖬𝖢𝖦}(,L)`$ is the quotient $`\mathrm{𝖧𝗈𝗆𝖾𝗈}_+(,L)/`$, where $`\mathrm{𝖧𝗈𝗆𝖾𝗈}_+(,L)`$ denotes the group of orientation-preserving homeomorphisms fixing set-wise $`L`$ and $``$ denotes homotopy. If we denote by $`\mathrm{𝖵𝖺𝗎𝗍}(\pi _1\mathrm{\Sigma })`$ the group of virtual automorphisms of $`\pi _1\mathrm{\Sigma }`$ (c.f. §10) then
###### Theorem.
There is an isomorphism
$$\mathrm{\Theta }:\mathrm{𝖬𝖢𝖦}(,L)\mathrm{𝖵𝖺𝗎𝗍}(\pi _1\mathrm{\Sigma }).$$
A proof of this theorem first appeared in the unpublished 1997 thesis of C. Odden . Due to its importance in the genus-independent expression of the Ehrenpreis conjecture , we provide a proof in order to ensure its inclusion in the literature.
Acknowledgements: I would like to thank P. Makienko and A. Verjovsky with whom I enjoyed fruitful conversations regarding several important aspects of this paper. I would also like to thank the Instituto de Matemáticas of the UNAM for providing a pleasant work enviroment and generous financial support.
## 2. Virtual Geometry
Virtual geometry is obtained as a quotient of nonstandard geometry, which we now review: references , , .
Let $`M`$ be a topological space, $`𝔘\mathrm{𝟤}^{}`$ an ultrafilter on the natural numbers all of whose elements have infinite cardinality. The nonstandard space $`{}_{}{}^{}M`$ is the set of sequences in $`M`$ modulo $`𝔘`$: that is,
(1)
$$\{x_i\}\{y_i\}\text{ if and only if }\{x_i\}|_X=\{y_i\}|_X\text{ for some }X𝔘.$$
Elements of $`{}_{}{}^{}M`$ are denoted $`{}_{}{}^{}x`$. There is a natural map $`M{}_{}{}^{}M`$ given by the constant sequences. Modulo the continuum hypothesis, $`{}_{}{}^{}M`$ is independent of the choice of ultrafilter.
There are two topologies on $`{}_{}{}^{}M`$ that naturally suggest themselves. The enlargement topology is generated by sets of the form $`{}_{}{}^{}O`$, where $`O`$ is open in $`M`$. It has the same countability as the topology of $`M`$ but is non-Hausdorff. The internal topology is generated by sets of the form $`[O_\alpha ]=\{{}_{}{}^{}x{}_{}{}^{}M|{}_{}{}^{}x\text{ is represented by a sequence }\{x_\alpha \},x_\alpha O_\alpha \}`$, where $`\{O_\alpha \}`$ is any sequence of open sets of $`M`$. It is Hausdorff but has greater countability than the topology of $`M`$.
For example, if we let $`M=`$ we obtain the nonstandard reals $`{}_{}{}^{}`$, a totally ordered, non-archemidean field. Note that $`{}_{}{}^{}`$ is an infinite-dimensional vector space over $``$. We will refer to the following substructures of the nonstandard reals:
* The subring of bounded nonstandard reals, denoted $`{}_{}{}^{}_{\mathrm{𝖿𝗂𝗇}}^{}`$, which consists of all classes of sequences that are bounded.
* The additive subgroup of infinitesimals, denoted $`{}_{}{}^{}_{ϵ}^{}`$, which consists of all classes of sequences converging to $`0`$.
* The cone of positive elements, denoted $`{}_{}{}^{}_{+}^{}`$, which consists of all classes of sequences that are $`0`$.
$`{}_{}{}^{}_{\mathrm{𝖿𝗂𝗇}}^{}`$ is a local topological ring in either the enlargement or internal topology, with maximal ideal $`{}_{}{}^{}_{ϵ}^{}`$. The quotient $`{}_{}{}^{}_{\mathrm{𝖿𝗂𝗇}}^{}/{}_{}{}^{}_{ϵ}^{}`$ is isomorphic to $``$, homeomorphic with the quotient enlargement topology (the quotient internal topology is discrete). The inclusion $`{}_{}{}^{}_{\mathrm{𝖿𝗂𝗇}}^{}`$ allows us to canonically identify $`{}_{}{}^{}_{\mathrm{𝖿𝗂𝗇}}^{}`$ with the product $`\times {}_{}{}^{}_{ϵ}^{}`$. Taking the product of the euclidean topology on $``$ with the discrete topology on $`{}_{}{}^{}_{ϵ}^{}`$, we obtain a third topology on $`{}_{}{}^{}_{\mathrm{𝖿𝗂𝗇}}^{}`$ which is Hausdorff and quotients by $`{}_{}{}^{}_{ϵ}^{}`$ to the topology on $``$. We call this third topology the lamination topology: it may be extended to $`{}_{}{}^{}`$ by giving the group $`{}_{}{}^{}/`$ the discrete topology and identifying $`{}_{}{}^{}\times ({}_{}{}^{}/)`$.
If $`M`$ is an $`n`$-manifold, then $`{}_{}{}^{}M`$ is a nonstandard manifold modelled on $`{}_{}{}^{}_{}^{n}`$. If we denote by $`{}_{}{}^{}M_{\mathrm{𝖿𝗂𝗇}}^{}`$ the points of $`{}_{}{}^{}M`$ represented by sequences which converge to points of $`M`$, then we may choose an atlas on $`{}_{}{}^{}M_{\mathrm{𝖿𝗂𝗇}}^{}`$ whose transitions preserve the lamination structure of $`{}_{}{}^{}_{\mathrm{𝖿𝗂𝗇}}^{n}`$ i.e. $`{}_{}{}^{}M_{\mathrm{𝖿𝗂𝗇}}^{}`$ is an $`n`$-lamination. In general, $`{}_{}{}^{}M`$ is a union of laminations of dimensions $`n`$, this because of the possibility of “dimension collapse” which we describe in the proof of Theorem 1 below.
If $`d`$ is a metric inducing the topology of $`M`$, it extends to a $`{}_{}{}^{}_{+}^{}`$-valued metric $`{}_{}{}^{}d`$ on $`{}_{}{}^{}M`$. Write $`{}_{}{}^{}x{}_{}{}^{}x_{}^{}`$ if $`{}_{}{}^{}d({}_{}{}^{}x,{}_{}{}^{}x_{}^{}){}_{}{}^{}_{ϵ}^{}`$.
###### Definition 1.
The virtual extension of $`M`$ is the quotient
$${}_{}{}^{}M={}_{}{}^{}M/,$$
equipped with the quotient lamination topology.
The virtual extension of $`{}_{}{}^{}`$ of $``$ is called the virtual reals, a totally-ordered real vector space. The metric $`{}_{}{}^{}d`$ on $`{}_{}{}^{}M`$ induces a $`{}_{}{}^{}_{+}^{}`$-valued metric $`{}_{}{}^{}d`$ on $`{}_{}{}^{}M`$. Given $`{}_{}{}^{}x{}_{}{}^{}M`$, the set
$$U_{}_{}{}^{}x=\{{}_{}{}^{}y|{}_{}{}^{}d({}_{}{}^{}x,{}_{}{}^{}y)\}$$
is a component of $`{}_{}{}^{}M`$ called the galaxy of $`{}_{}{}^{}x`$. $`M`$ is a galaxy of $`{}_{}{}^{}M`$, and $`{}_{}{}^{}M`$ is the union of all of its galaxies.
The galaxies of $`{}_{}{}^{}M`$ can be quite different from one another. For example if $`M`$ is simply connected, then there may be galaxies that are not. For example, suppose that $`M`$ is a noncompact leaf of the Reeb foliation of the torus. Consider a sequence of points $`\{x_\alpha \}`$ in $`M`$ converging to a point $`\widehat{x}`$ in the compact toral leaf. Let $`\{\gamma _\alpha \}`$ be a sequence of simple closed curves converging to the meridian through $`\widehat{x}`$. Then the limit curve $`{}_{}{}^{}\gamma `$ is essential in the universe $`U_{}_{}{}^{}x`$. On the other hand, if $`M`$ is a riemannian homogeneous space, then the universes of $`{}_{}{}^{}M`$ are all isometric to $`M`$.
###### Theorem 1.
If $`M`$ is a complete riemannian manifold of dimension $`n`$, each galaxy $`U`$ of $`{}_{}{}^{}M`$ has the structure of a complete riemannian manifold of dimension $`mn`$.
###### Proof.
Given a galaxy $`U`$, $`{}_{}{}^{}xU`$ and $`\{x_\alpha \}`$ a representative sequence, let $`m`$ be the largest integer for which there exists a sequence of $`m`$-dimensional balls $`\{D_r(x_\alpha )\}`$ of fixed radius $`r`$ about $`\{x_\alpha \}`$. The integer $`m`$ is independent of the representative sequence and defines an $`m`$-ball $`D_r({}_{}{}^{}x)U`$. The function $`{}_{}{}^{}xm`$ is locally constant, thus the collection of such balls defines on $`U`$ the structure of a smooth $`m`$-manifold. Note that it is possible to have $`m<n`$: for example, if $`M`$ is a hyperbolic manifold with a cusp, then for a class of sequence emptying into the cusp, we have $`m=n1`$.
Consider the nonstandard tangent bundle
$$𝐓{}_{}{}^{}M:={}_{}{}^{}\left(𝐓M\right).$$
There is a natural projection of $`𝐓{}_{}{}^{}M`$ onto $`{}_{}{}^{}M`$ whose fiber $`𝐓_{}_{}{}^{}x{}_{}{}^{}M`$ – the tangent space at $`{}_{}{}^{}x`$ – consists of classes of sequences of vectors $`\{𝗏_\alpha \}`$ based at sequences $`\{x_\alpha \}`$ belonging to the class of $`{}_{}{}^{}x`$. It is not difficult to see that $`𝐓_{}_{}{}^{}x{}_{}{}^{}M`$ is a real infinite-dimensional vector space. The riemannian metric $`\rho `$ extends to a $`{}_{}{}^{}`$-valued metric $`{}_{}{}^{}\rho `$ on $`𝐓{}_{}{}^{}M`$ in the obvious way. Denote by $`{}_{}{}^{}||`$ the associated norm. Define the bounded tangent bundle by
$$𝐓_{\mathrm{𝖿𝗂𝗇}}{}_{}{}^{}M=\{{}_{}{}^{}𝗏𝐓{}_{}{}^{}M|{}_{}{}^{}|{}_{}{}^{}𝗏|{}_{}{}^{}_{\mathrm{𝖿𝗂𝗇}}^{}\}.$$
Given tangent vectors $`{}_{}{}^{}𝗏`$ and $`{}_{}{}^{}𝗏_{}^{}`$ based at $`{}_{}{}^{}x`$ and $`{}_{}{}^{}x_{}^{}`$, we write $`{}_{}{}^{}𝗏{}_{}{}^{}𝗏_{}^{}`$ if
1. $`{}_{}{}^{}x{}_{}{}^{}x_{}^{}`$.
2. the Levi-Civita parallel translate of a representative $`\{𝗏_\alpha \}`$ of $`{}_{}{}^{}𝗏`$ to a representative $`\{x_\alpha ^{}\}`$ of $`{}_{}{}^{}x_{}^{}`$ – along a sequence of geodesics connecting to a representative $`\{x_\alpha \}`$ of $`{}_{}{}^{}x`$ – is asymptotic to a representative $`\{𝗏_\alpha ^{}\}`$ of $`{}_{}{}^{}𝗏_{}^{}`$.
Now define the bounded tangent bundle of $`{}_{}{}^{}M`$ to be
$$𝐓_{\mathrm{𝖿𝗂𝗇}}{}_{}{}^{}M=𝐓_{\mathrm{𝖿𝗂𝗇}}{}_{}{}^{}M/.$$
The nonstandard riemannian metric $`{}_{}{}^{}\rho `$ on $`𝐓_{\mathrm{𝖿𝗂𝗇}}{}_{}{}^{}M`$ descends to a riemannian metric on $`𝐓_{\mathrm{𝖿𝗂𝗇}}{}_{}{}^{}M`$. If $`U`$ is a galaxy, its tangent space may be identified with the restriction of $`𝐓_{\mathrm{𝖿𝗂𝗇}}{}_{}{}^{}M`$ to $`U`$. Now any geodesic $`\eta U`$ can be realized as a sequence class of geodesics $`\{\eta _\alpha \}`$. Since each member of such a sequence can be continued indefinitely, the same is true of $`\eta `$, hence $`U`$ is complete. ∎
###### Definition 2.
Let $`M`$, $`N`$ be riemannian $`n`$-manifolds. A virtual subisometry is an injective map
$${}_{}{}^{}f:{}_{}{}^{}M{}_{}{}^{}N,$$
where $`{}_{}{}^{}f`$ maps each galaxy of $`{}_{}{}^{}M`$ isometrically onto a galaxy of $`{}_{}{}^{}N`$. If in addition there exists a virtual subisometry $`{}_{}{}^{}g:{}_{}{}^{}N{}_{}{}^{}M`$, then the pair $`({}_{}{}^{}f,{}_{}{}^{}g)`$ is called a virtual isometry.
We write $`M_{\mathrm{𝗏𝗂𝗋}}N`$ to indicate the existance of a virtual subisometry $`{}_{}{}^{}f`$ and $`M_{\mathrm{𝗏𝗂𝗋}}N`$ indicates the existence of a virtual isometry. The relation $`_{\mathrm{𝗏𝗂𝗋}}`$ defines a partial ordering on the set of all riemannian $`n`$-manifolds.
An isometry $`f:MN`$ clearly induces a virtual isometry $`({}_{}{}^{}f,{}_{}{}^{}g):{}_{}{}^{}M{}_{}{}^{}N`$ with $`{}_{}{}^{}f,{}_{}{}^{}g`$ inverse to one another. More generally, a continuous map $`{}_{}{}^{}f:{}_{}{}^{}M{}_{}{}^{}N`$ is called standard if it is induced by a map $`f:MN`$ i.e. if for any $`{}_{}{}^{}x{}_{}{}^{}M`$ and any representative $`\{x_\alpha \}`$, $`\{f(x_\alpha )\}`$ is a representative of $`{}_{}{}^{}f({}_{}{}^{}x)`$.
###### Theorem 2.
Let $`L`$ be a dense leaf of a riemannian lamination $``$. Then for every leaf $`L^{}`$,
$$\stackrel{~}{L}^{}_{\mathrm{𝗏𝗂𝗋}}\stackrel{~}{L}.$$
###### Proof.
Fix a global metric $`d`$ on $``$ which agrees locally with the riemannian metric on the leaves. (By this we mean that in sufficiently small flow boxes, $`d`$ agrees with the distance function of $`\rho `$ in any plaque.) Let $`\{\stackrel{~}{x}_\alpha ^{}\}\stackrel{~}{L}^{}`$ be any sequence, $`\{x_\alpha ^{}\}`$ its projection to $`L^{}`$. Let $`\{\stackrel{~}{x}_\alpha \}\stackrel{~}{L}`$ be a sequence whose projection $`\{x_\alpha \}`$ to $`L`$ is $`d`$-asymptotic to $`\{x_\alpha ^{}\}`$. By transversal continuity of the metric, we deduce a sequence of $`K_\alpha `$-quasiisometries, $`K_\alpha 1`$,
$$f_\alpha :D_\delta (x_\alpha )D_\delta (x_\alpha ^{}),$$
for some $`\delta >0`$, where $`D_\delta (x)`$ means the open $`\rho `$-ball of radius $`\delta `$ about $`x`$. Then if $`{}_{}{}^{}\stackrel{~}{x}`$, $`{}_{}{}^{}\stackrel{~}{x}_{}^{}`$ are the virtual classes of $`\{\stackrel{~}{x}_\alpha \}`$, $`\{\stackrel{~}{x}_\alpha ^{}\}`$, the sequence of quasiisometries $`\{f_\alpha \}`$ induces an isometry $`D_\delta ({}_{}{}^{}\stackrel{~}{x})D_\delta ({}_{}{}^{}\stackrel{~}{x}_{}^{})`$. Since $`L`$ is dense, we may continue these isometries along geodesics to obtain a locally isometric surjection $`UU^{}`$, where $`U`$, $`U^{}`$ are the galaxies containing $`{}_{}{}^{}\stackrel{~}{x}`$, $`{}_{}{}^{}\stackrel{~}{x}_{}^{}`$. But since these spaces are simply connected, and the map is isometric, this surjection is a bijection. Hence it inverts to an isometry $`U^{}U`$. Repeating this for every $`^{}`$-class of sequence in $`\stackrel{~}{L}^{}`$, we obtain the desired virtual subisometry $`\stackrel{~}{L}^{}_{\mathrm{𝗏𝗂𝗋}}\stackrel{~}{L}`$. ∎
Two riemannian manifolds have the same virtual geometry if their universal covers are virtually isometric.
###### Corollary 1.
Dense leaves of a riemannian lamination $``$ have the same virtual geometry.
## 3. The Fundamental Germ
Let $``$ be a riemannian lamination, $`x`$ a point contained in a transversal $`T`$, $`L`$ a leaf accumulating at $`x`$ and $`L_x`$ the leaf containing $`x`$. Let $`𝖿:T𝐅_{}`$ be a continuous section of the leaf-wise orthonormal frame bundle of $``$ over $`T`$. Fix locally isometric universal covers $`p:\stackrel{~}{L}L`$ and $`p_x:\stackrel{~}{L}_xL_x`$. Denote $`T_0=TL`$, $`\stackrel{~}{T}_0=p^1(T_0)`$ and let $`\stackrel{~}{𝖿}_{\stackrel{~}{y}}`$ denote the lift of $`𝖿_y`$ to a point $`\stackrel{~}{y}\stackrel{~}{T}_0`$ covering $`y`$. We pick a basepoint $`\stackrel{~}{x}\stackrel{~}{L}_x`$ lying over $`x`$ with lifted frame $`\stackrel{~}{𝖿}_{\stackrel{~}{x}}`$.
Let $`\stackrel{~}{y}\stackrel{~}{T}_0`$. For $`r>0`$, the frames $`\stackrel{~}{𝖿}_{\stackrel{~}{x}}`$, $`\stackrel{~}{𝖿}_{\stackrel{~}{y}}`$ determine polar coordinates on the metric disks $`D_r(\stackrel{~}{x})`$, $`D_r(\stackrel{~}{y})`$. This yields in turn a canonical quasiisometry
$$f:D_r(\stackrel{~}{x})D_r(\stackrel{~}{y})$$
given by the coordinate maps.
Let $`\{x_\alpha \}T_0`$ be a sequence converging to $`x`$, $`\{\stackrel{~}{x}_\alpha \}\stackrel{~}{T}_0`$ any sequence covering $`\{x_\alpha \}`$. Then the frame sequence $`\{\stackrel{~}{𝖿}_{\stackrel{~}{x}_\alpha }\}`$ and the frame $`\stackrel{~}{f}_{\stackrel{~}{x}}`$ determine a sequence of $`K_\alpha `$-quasiisometries
$$\{f_\alpha :D_{r_\alpha }(\stackrel{~}{x})D_{r_\alpha }(\stackrel{~}{x}_\alpha )\}.$$
Since $`L`$ accumulates at $`x`$, we may choose the sequence of radii $`r_\alpha \mathrm{}`$ so that $`K_\alpha 1`$. We deduce an isometry
$${}_{}{}^{}f:\stackrel{~}{L}_xU{}_{}{}^{}\stackrel{~}{L}$$
where $`U`$ is the galaxy containing $`{}_{}{}^{}\stackrel{~}{x}`$. The map $`{}_{}{}^{}f`$ is called an f-diophantine approximation of $`x`$ along $`L`$.
###### Definition 3.
The fundamental germ of $``$, based at $`x`$ along $`L`$ and $`𝖿`$, is
$$[[\pi ]]_1(,L,x,𝖿)=\left\{{}_{}{}^{}f{}_{}{}^{}g_{}^{1}\right|{}_{}{}^{}f,{}_{}{}^{}g\text{ are }𝖿\text{-diophantine approximations of }x\text{ along }L\}.$$
If $`xL`$, we shorten the notation to $`[[\pi ]]_1(,x,𝖿)`$. The groupoid structure of $`[[\pi ]]_1(,L,x,𝖿)`$ will be described in the next section.
###### Note 1.
Suppose that $``$ is a constant curvature riemannian foliation with dense leaf $`L`$ modeled on the space form $`𝕄^n=^n`$ or $`^n`$. Then the frame field actually determines a sequence of uniquely defined global isometries $`\{f_\alpha :𝕄^n𝕄^n\}`$. Given $`G`$ a group, nonstandard $`G`$ is the group $`{}_{}{}^{}G`$ of all sequences $`\{g_\alpha \}G`$ modulo the relation $``$ described in (1). Then an $`𝖿`$-diophantine approximation is completely determined by the class $`{}_{}{}^{}f{}_{}{}^{}\mathrm{𝖨𝗌𝗈𝗆}(𝕄^n)`$ of $`\{f_\alpha \}`$. We note that $`{}_{}{}^{}\mathrm{𝖨𝗌𝗈𝗆}(𝕄^n)`$ is a subgroup of $`\mathrm{𝖨𝗌𝗈𝗆}({}_{}{}^{}𝕄_{}^{n})`$ (the group of isometries of $`{}_{}{}^{}𝕄_{}^{n}`$, not virtual isometries). Thus, if $`\mathrm{\Gamma }\pi _1L`$ is the deck group of $`𝕄^nL`$, we have
$${}_{}{}^{}\mathrm{\Gamma }[[\pi ]]_1(,L,x,𝖿){}_{}{}^{}\mathrm{𝖨𝗌𝗈𝗆}(𝕄^n).$$
The terminology $`𝖿`$-diophantine approximation comes from the following example.
###### Example 1.
Let $``$ be the irrational foliation of the torus $`𝕋^2`$ by lines of slope $`r`$. Define a representation $`\rho :\pi _1𝕊^1\mathrm{𝖧𝗈𝗆𝖾𝗈}(𝕊^1)`$ by $`\rho _n(\overline{y})=\overline{ynr}`$, where $`\overline{y}`$ denotes the image of $`y`$ in $`𝕊^1=/`$. Then the suspension of $`\rho `$, $`_\rho =(\times 𝕊^1)/`$ , is homeomorphic to $``$. The map $`\times 𝕊^1𝕊^1`$ defined $`(x,\overline{y})\overline{x}`$ (i.e. the projection onto the first factor composed with the universal covering $`𝕊^1`$) induces a projection $`_\rho 𝕊^1`$. Let $`T𝕊^1`$ be a fiber of this projection passing through $`x`$. A frame section $`𝖿`$ along $`T`$ is determined by an orientation of $``$. In this case, an $`𝖿`$-diophantine approximation of $`x`$ is just a diophantine approximation of $`r`$. (Recall that a sequence $`\{n_\alpha \}`$ is called a diophantine approximation of $`r`$ if $`\{\overline{rn_\alpha }\}`$ converges to $`\overline{0}𝕊^1`$.) Thus if one denotes by $`{}_{}{}^{}_{r}^{}`$ the subgroup of $`{}_{}{}^{}`$ consisting of classes of diophantine approximations of $`r`$, we obtain in agreement with the construction in , §4.4:
$$[[\pi ]]_1(,L,x,𝖿)={}_{}{}^{}_{r}^{}.$$
(Note: $`{}_{}{}^{}_{r}^{}`$ is an ideal if and only if $`r`$ is rational.) If another frame field $`𝖿^{}`$ is used whose domain is a transversal $`T^{}`$ which is not a suspension fiber, the set of diophantine approximations is a subset $`{}_{}{}^{}_{r}^{}{}_{}{}^{}`$. This subset maps injectively into $`{}_{}{}^{}`$ with image $`{}_{}{}^{}_{r}^{}={}_{}{}^{}_{r}^{}`$.
###### Example 2.
Consider a nested set of Fuchsian groups $`𝒢=\{\mathrm{\Gamma }_i\}`$ and let
$$\widehat{\mathrm{\Sigma }}_𝒢=\underset{}{lim}^2/\mathrm{\Gamma }_i,$$
be the associated hyperbolic surface solenoid. We may take $`T`$ to be a fiber $`\widehat{p}^1(x_0)`$ of the projection $`\widehat{p}:\widehat{\mathrm{\Sigma }}_𝒢\mathrm{\Sigma }_0`$, where $`\mathrm{\Sigma }_0=^2/\mathrm{\Gamma }_0`$ is the initial surface. Then a frame at $`x_0`$ pulls back to a frame section $`𝖿`$ along $`T`$. In this case, we find that Definition 3 again agrees with the definition found in :
$`[[\pi ]]_1(\widehat{\mathrm{\Sigma }}_𝒢,L,x,𝖿)`$ $`=`$ $`{\displaystyle {}_{}{}^{}\mathrm{\Gamma }_{i}^{}}`$
$`=`$ $`\{\{g_\alpha \}\mathrm{\Gamma }_0|\text{for all }i,N_i\text{ such that }g_\alpha \mathrm{\Gamma }_i\text{ when }\alpha >N_i\}/,`$
a subgroup of $`{}_{}{}^{}PSL(2,)PSL(2,{}_{}{}^{})`$. If $`𝖿^{}`$ is another frame field, not necessarily with a fiber transversal domain, then the corresponding germ $`[[\pi ]]_1(\widehat{\mathrm{\Sigma }}_𝒢,L,x,𝖿^{})`$ need not define a subgroup of $`PSL(2,{}_{}{}^{})`$ and particularly, need not be isomorphic to $`{}_{}{}^{}\mathrm{\Gamma }_{i}^{}`$ (although the fundamental germs calculated with respect to $`𝖿`$ and $`𝖿^{}`$ are in canonical bijection). The rub here is the non-uniform nature of the action of $`PSL(2,)`$ on $`^2`$. This problem will become moot through the replacement of the fundamental germ by the mother germ, §6.
###### Example 3.
More generally, let $``$ be any hyperbolic surface lamination. Then the fundamental germ $`[[\pi ]]_1(,L,x,𝖿)`$ is a subset of $`PSL(2,{}_{}{}^{})`$. Equally, if $``$ is a hyperbolic 3-lamination, $`[[\pi ]]_1(,L,x,𝖿)PSL(2,{}_{}{}^{})`$.
## 4. The Germ Universal Cover
Let $`L`$ be a fixed leaf. Denote by $`p:\stackrel{~}{L}L`$ the universal cover. We recall the following definition :
###### Definition 4.
The germ universal cover of $``$ along $`L`$ is the subspace $`[[\stackrel{~}{}]]{}_{}{}^{}\stackrel{~}{L}`$ defined
$$[[\stackrel{~}{}]]=\{\{\stackrel{~}{x}_\alpha \}\stackrel{~}{L}|\left\{p(\stackrel{~}{x}_\alpha )\right\}\text{ converges in }\}/.$$
We will denote elements of the germ universal cover by $`{}_{}{}^{}\stackrel{~}{x}`$. There is a natural projection
$${}_{}{}^{}p:[[\stackrel{~}{}]]{}_{}{}^{}\stackrel{~}{x}\widehat{x}=limp(\stackrel{~}{x}_\alpha ),$$
where $`\{\stackrel{~}{x}_\alpha \}`$ is a representative sequence in the class $`{}_{}{}^{}\stackrel{~}{x}`$. We will write $`lim{}_{}{}^{}\stackrel{~}{x}=\widehat{x}`$ if $`{}_{}{}^{}p({}_{}{}^{}\stackrel{~}{x})=\widehat{x}`$. Note that $`{}_{}{}^{}p`$ is surjective if and only if $`L`$ is dense and in general $`{}_{}{}^{}p`$ maps onto the closure $`\overline{L}`$ of $`L`$, itself a sublamination of $``$.
###### Proposition 1.
$`[[\stackrel{~}{}]]`$ consists of a union of galaxies of $`{}_{}{}^{}\stackrel{~}{L}`$.
###### Proof.
Let $`{}_{}{}^{}\stackrel{~}{x}[[\stackrel{~}{}]]`$ and denote by $`U`$ the galaxy containing $`{}_{}{}^{}\stackrel{~}{x}`$. If $`{}_{}{}^{}\stackrel{~}{y}U`$, then there exists a sequence of geodesic paths $`\{\stackrel{~}{\eta }_\alpha \}`$ connecting representatives $`\{\stackrel{~}{x}_\alpha \}`$ to $`\{\stackrel{~}{y}_\alpha \}`$ in $`\stackrel{~}{L}`$, whose projection to $``$ gives a convergent sequence of paths $`\{\eta _\alpha \}`$. It follows that the projection $`\{p(\stackrel{~}{y}_\alpha )\}`$ converges, and $`{}_{}{}^{}\stackrel{~}{y}[[\stackrel{~}{}]]`$ as well. ∎
The galaxies that make up $`[[\stackrel{~}{}]]`$ will be referred to as leaves. See for a proof of the following
###### Theorem 3.
$`[[\stackrel{~}{}]]`$ may be given the structure of a lamination whose leaves are nowhere dense and such that the map $`{}_{}{}^{}p:[[\stackrel{~}{}]]\overline{L}`$ is an open surjection.
One can thus think of $`[[\stackrel{~}{}]]`$ as a the result of unwrapping all of the diophantine approximations implied by $`L`$. The topology that $`[[\stackrel{~}{}]]`$ obtains from its lamination atlas is not unique, and is called a germ universal cover topology. It is in general coarser than the topology $`[[\stackrel{~}{}]]`$ induces from $`{}_{}{}^{}\stackrel{~}{L}`$.
###### Proposition 2.
If $``$ is compact then $`[[\stackrel{~}{}]]={}_{}{}^{}\stackrel{~}{L}`$.
###### Proof.
This follows from well-known compactness arguments e.g. see the proof in . ∎
An element $`{}_{}{}^{}u={}_{}{}^{}f{}_{}{}^{}g_{}^{1}[[\pi ]]_1(,L,x,𝖿)`$ arises as the limit of a sequence of $`K_\alpha `$-quasiisometries
(2)
$$\{u_\alpha :D_{r_\alpha }(\stackrel{~}{x}_\alpha )D_{r_\alpha }(\stackrel{~}{y}_\alpha )\},$$
where $`\{\stackrel{~}{x}_\alpha \},\{\stackrel{~}{y}_\alpha \}L`$, $`K_\alpha 1`$ and $`r_\alpha \mathrm{}`$. The limit $`{}_{}{}^{}u:VU`$ is independent of the sequence $`\{u_\alpha \}`$ and depends only on the sequences of frames $`\{𝖿_{x_\alpha }\}`$, $`\{𝖿_{y_\alpha }\}`$. In particular, we could have obtained $`{}_{}{}^{}u`$ through the same sequence of quasiisometries with domains extended to a sequence of larger disks $`D_{s_\alpha }(\stackrel{~}{x}_\alpha )`$, $`s_\alpha >r_\alpha `$ – provided that the new quasiisometry constants converge to 1 as well.
Now for arbitrary $`{}_{}{}^{}\stackrel{~}{x}[[\stackrel{~}{}]]`$, the expression $`{}_{}{}^{}u({}_{}{}^{}x)`$ does not even make formal sense, since $`{}_{}{}^{}u`$ is so far only defined on the galaxy $`V`$. We contrast this with the constant curvature case, where, because $`[[\pi ]]_1(,x,L,𝖿)\mathrm{𝖨𝗌𝗈𝗆}({}_{}{}^{}𝕄)`$, $`{}_{}{}^{}u({}_{}{}^{}\stackrel{~}{x})`$ is always formally defined, although it need not define an element of $`[[\stackrel{~}{}]]`$.
Let us say that $`{}_{}{}^{}u`$ is formally defined on an element $`{}_{}{}^{}\stackrel{~}{w}[[\stackrel{~}{}]]`$ if there exists a sequence (2) giving rise to $`{}_{}{}^{}u`$ and a representative sequence $`\{\stackrel{~}{w}_\alpha \}`$ of $`{}_{}{}^{}\stackrel{~}{w}`$ such that
$$D_{r_\alpha ^{}}(w_\alpha )D_{r_\alpha }(\stackrel{~}{x}_\alpha )$$
for all $`\alpha `$, where $`r_\alpha ^{}\mathrm{}`$. It follows then that if $`V^{}`$ is the galaxy containing $`{}_{}{}^{}\stackrel{~}{w}`$, then the limit $`{}_{}{}^{}u`$ is defined on $`V^{}`$ as well. Whenever we write $`{}_{}{}^{}u({}_{}{}^{}w)`$, it will tacitly be understood that $`{}_{}{}^{}u`$ is formally defined at $`{}_{}{}^{}w`$.
Define the domain of $`{}_{}{}^{}u`$ as
$$\mathrm{𝖣𝗈𝗆}({}_{}{}^{}u)=\left\{{}_{}{}^{}\stackrel{~}{x}[[\stackrel{~}{}]]\right|{}_{}{}^{}u({}_{}{}^{}\stackrel{~}{x})[[\stackrel{~}{}]]\text{ and }lim{}_{}{}^{}u({}_{}{}^{}\stackrel{~}{x})=lim{}_{}{}^{}\stackrel{~}{x}\},$$
and $`\mathrm{𝖱𝖺𝗇}({}_{}{}^{}u)={}_{}{}^{}u(\mathrm{𝖣𝗈𝗆}({}_{}{}^{}u))`$. With this definition, it follows that $`[[\pi ]]_1(,L,x,𝖿)`$ has the structure of a groupoid. Note that for any $`{}_{}{}^{}u[[\pi ]]_1(,L,x,𝖿)`$, $`\mathrm{𝖣𝗈𝗆}({}_{}{}^{}u)`$, $`\mathrm{𝖱𝖺𝗇}({}_{}{}^{}u)`$ are unions of leaves and hence induce lamination structures from $`[[\stackrel{~}{}]]`$. Moreover, on $`\mathrm{𝖣𝗈𝗆}({}_{}{}^{}u)`$,
(3)
$${}_{}{}^{}p{}_{}{}^{}u={}_{}{}^{}p.$$
In particular we see that $`{}_{}{}^{}u:\mathrm{𝖣𝗈𝗆}({}_{}{}^{}u)\mathrm{𝖱𝖺𝗇}({}_{}{}^{}u)`$ defines a lamination homeomorphism.
###### Example 4.
Let $``$ be the irrational foliation of $`𝕋^2`$ by lines of slope $`r`$, $`L`$ any dense leaf. Then by Proposition 2, $`[[\stackrel{~}{}]]={}_{}{}^{}`$. Moreover, for any frame field $`𝖿`$ and $`{}_{}{}^{}u[[\pi ]]_1(,L,x,𝖿)`$, it is not difficult to see that $`\mathrm{𝖣𝗈𝗆}({}_{}{}^{}u)={}_{}{}^{}`$. Thus, $`[[\pi ]]_1(,L,x,𝖿)`$ is a group isomorphic to $`{}_{}{}^{}_{r}^{}`$.
###### Example 5.
Let $``$ be the profinite hyperbolic surface solenoid $`\widehat{\mathrm{\Sigma }}_𝒢`$ of Example 2. Then we have, again by compactness,
$$[[\stackrel{~}{}]]={}_{}{}^{}_{}^{2}.$$
If $`𝖿`$ is a frame field lifted from a frame on a surface occurring in the defining inverse limit, then $`[[\pi ]]_1(,L,x,𝖿)`$ is a group. On the other hand, if $`𝖿`$ is a frame field not obtained in this way, then $`[[\pi ]]_1(,L,x,𝖿)`$ need not be a group e.g. see §6.
The following may also be found in .
###### Theorem 4.
Let $`F:(,L)(^{},L^{})`$ be a lamination map. Then there exist germ universal cover topologies so that the map
$$[[\stackrel{~}{F}]]:[[\stackrel{~}{}]][[\stackrel{~}{^{}}]]$$
induced by $`\{\stackrel{~}{x}_\alpha \}\{\stackrel{~}{F}(\stackrel{~}{x}_\alpha )\}`$ is a continuous lamination map.
###### Note 2.
It is useful here to point out that for a lamination $`_\rho =(\stackrel{~}{B}\times 𝖥)/\pi _1B`$ occurring as a suspension of a representation $`\rho :\pi _1B\mathrm{𝖧𝗈𝗆𝖾𝗈}(𝖥)`$, it is in general false that a lamination homeomorphism $`F:_\rho _\rho `$ lifts to a homeomorphism of the “universal covering space” $`\stackrel{~}{B}\times 𝖥`$.
Now suppose that $`L^{}`$ is another leaf of $``$. Denote by $`[[\stackrel{~}{}]]^{}`$ the germ universal cover formed from $`L^{}`$.
###### Proposition 3.
If $`L^{}`$ accumulates on $`L`$ then there is a virtual subisometry $`{}_{}{}^{}\stackrel{~}{L}{}_{}{}^{}\stackrel{~}{L}_{}^{}`$ restricting to a virtual subisometry
$$[[\stackrel{~}{}]][[\stackrel{~}{}]]^{}$$
which is a homeomorphism onto its image with respect to appropriate germ universal cover topologies.
###### Proof.
This follows directly from the proof of Theorem 2. ∎
## 5. Sensitivity to Changes in Data
In this section we shall examine the dependence of the fundamental germ on the base point $`x`$, the accumulating leaf $`L`$ and the frame field $`𝖿`$.
Change in base point and accumulating leaf: Let us fix for the moment the dense leaf $`L`$ and consider a change of base point $`xx^{}`$ in which $`L_x=L_x^{}`$. Let $`\eta `$ be a geodesic connecting $`x`$ to $`x^{}`$ in $`L_x`$. The tangent vector $`𝗏`$ to $`\eta `$ at $`x`$ has coordinate $`(a_1,\mathrm{},a_n)`$ with respect to the frame $`𝖿_x`$. At each $`yT=`$ the domain of $`𝖿`$, this coordinate determines a vector $`𝗏_y`$ using the frame $`𝖿_y`$. We obtain in this way a transversally continuous family of geodesics $`\{\eta _y\}_{yT}`$. Restricting to an open subtransversal of $`T`$ if necessary, we may parallel translate $`𝖿`$ along the geodesic family to obtain a frame field $`𝖿^{}`$ with domain $`T^{}x^{}`$. The following is then immediate from the definition of the fundamental germ.
###### Proposition 4.
Let $`x^{}`$, $`𝖿^{}`$ be as in the preceding paragraph. Then
$$[[\pi ]]_1(,L,x,𝖿)=[[\pi ]]_1(,L,x^{},𝖿^{}).$$
If we consider a change of base point $`xx^{}`$, in which $`L_xL_x^{}`$, the situation becomes considerably more subtle. In fact, we shall see in §7 that fundamental germs based at points on different leaves can be nonisomorphic. For similar reasons, a change in accumulating leaf $`L`$ may yield nonisomorphic fundamental germs.
Change in frame field: Let us now fix the base point $`x`$ and consider a new frame field $`𝖿^{}:T^{}𝐅_{}`$ based at $`x`$. For simplicity, we again assume that $`\pi _1L=1`$. Since $`T`$ (the domain of $`𝖿`$) and $`T^{}`$ each contain subtransversal neighborhoods of $`x`$ lying in a common flow box, it is clear that there is a natural bijection
$$[[\pi ]]_1(,L,x,𝖿)[[\pi ]]_1(,L,x,𝖿^{}).$$
The issue is then the law of composition. We will show that this map need not be an isomorphism.
Let us consider the inverse limit solenoid $`\widehat{\mathrm{\Sigma }}_𝒢`$ of Example 2. Assume that $`L=L_x`$, $`T=T^{}=`$ a fiber over a point $`x_0\mathrm{\Sigma }_0`$ and that $`𝖿_x=𝖿_x^{}`$. We will take $`𝖿`$ to be simply the lift of a frame based at $`x_0`$, so that $`𝖿`$-diophantine approximations consist of sequences $`\{\gamma _\alpha \}\mathrm{\Gamma }_0`$ converging with respect to the family $`\{\mathrm{\Gamma }_i\}`$. It follows that every $`𝖿^{}`$-diophantine approximation of $`x`$ may be written in the form $`\left\{\gamma _\alpha \mathrm{\Theta }_\alpha \right\}`$, where $`\{\mathrm{\Theta }_\alpha \}`$ consists of a sequence of rotations based at $`x`$ with angle going to 0 and $`\{\gamma _\alpha \}`$ is an $`𝖿`$-diophantine approximation. General elements of $`[[\pi ]]_1(,L,x,𝖿^{})`$ are then of the form $`\left\{\gamma _\alpha \mathrm{\Theta }_\alpha \eta _\alpha \right\}`$, where $`\{\eta _\alpha \}`$ is another $`𝖿`$-diophantine approximation. We should not expect products of elements of this type to yield elements of $`[[\pi ]]_1(,L,x,𝖿^{})`$. Indeed, such a product would have the shape
(4)
$$\left\{\gamma _\alpha \mathrm{\Theta }_\alpha \eta _\alpha \mathrm{\Delta }_\alpha \omega _\alpha \right\},$$
for $`\{\mathrm{\Delta }_\alpha \}`$ another sequence of rotations based at $`x`$ with angle going to $`0`$ and $`\{\omega _\alpha \}`$ another $`𝖿`$-diophantine approximation. If $`\mathrm{\Theta }_\alpha `$ does not converge to the identity fast enough, $`\mathrm{\Theta }_\alpha \eta _\alpha \mathrm{\Delta }_\alpha \omega _\alpha `$ applied to $`𝖿_x^{}`$ will not project to a frame based at $`x_0\mathrm{\Sigma }_0`$. Hence the expression (4) is not even asymptotic to an element of $`[[\pi ]]_1(,L,x,𝖿^{})`$. It is not difficult to see that unless $`𝖿^{}`$ is the pull-back of a frame on $`\mathrm{\Sigma }_0`$, this sort of problem always arises.
## 6. The Mother Germ
In this section, we assume that $``$ has a dense leaf $`L`$, with which we define the germ universal cover $`[[\stackrel{~}{}]]`$, equipped with a fixed germ universal cover topology.
While the fundamental germ $`[[\pi ]]_1(,L,x,𝖿)`$ enjoys the property of being reasonably calculable and leaf specific, it can be sensitive to data variation. There are additional shortcomings:
* By (3), the action of the fundamental germ $`[[\pi ]]_1(,L,x,𝖿)`$ on $`[[\stackrel{~}{}]]`$ respects the germ covering $`{}_{}{}^{}p`$. However it need not be the case that every identification implied by $`{}_{}{}^{}p`$ is implemented by an element of $`[[\pi ]]_1(,L,x,𝖿)`$.
* There will be in general other maps of leaves of $`[[\stackrel{~}{}]]`$ that satisfy (3) but do not appear in $`[[\pi ]]_1(,L,x,𝖿)`$.
* It appears that $`[[\pi ]]_1(,L,x,𝖿)`$ such as it is defined, will be functorial only under certain types of lamination maps e.g. see .
For this reason, we will expand $`[[\pi ]]_1(,L,x,𝖿)`$ to a larger groupoid, called the mother germ. The mother germ will be the maximal amplification of $`[[\pi ]]_1(,L,x,𝖿)`$ which contains all partially defined maps of sublaminations of $`[[\stackrel{~}{}]]`$ satisfying (3): in other words, it is the full deck groupoid of $`{}_{}{}^{}p`$.
Let $`\mathrm{𝖣𝗈𝗆},\mathrm{𝖱𝖺𝗇}`$ be sublaminations of $`[[\stackrel{~}{}]]`$. A homeomorphism
$${}_{}{}^{}u:\mathrm{𝖣𝗈𝗆}\mathrm{𝖱𝖺𝗇}[[\stackrel{~}{}]]$$
satisfying (3) is called deck. Note that condition (3) implies that a deck homeomorphism $`{}_{}{}^{}u`$ is automatically an isometry along the leaves of $`\mathrm{𝖣𝗈𝗆}`$.
###### Definition 5.
The mother germ is the groupoid
$$[[\pi ]]_1()=\{{}_{}{}^{}u:\mathrm{𝖣𝗈𝗆}\mathrm{𝖱𝖺𝗇}\text{ is a deck homeomorphism }\}.$$
The mother germ will never be a group, since it distinguishes deck maps obtained from others by restriction of domain. In general, however, it will contain many interesting and calculable subgroups and subgroupoids, as the following shows.
###### Proposition 5.
Let $`L^{}`$ be any dense leaf of $``$. Then there is an injective groupoid homomorphism
$$[[\pi ]]_1(,L^{},x,𝖿)[[\pi ]]_1().$$
###### Proof.
By Proposition 3, there exists an isometric inclusion $`{}_{}{}^{}f:[[\stackrel{~}{}]]^{}[[\stackrel{~}{}]]`$. If $`{}_{}{}^{}u[[\pi ]]_1(,L^{},x,𝖿)`$, then the map
$${}_{}{}^{}u{}_{}{}^{}f{}_{}{}^{}u{}_{}{}^{}f_{}^{1}$$
defines an injective groupoid homomorphism. ∎
###### Theorem 5.
The quotient
$$[[\pi ]]_1()\backslash [[\stackrel{~}{}]],$$
equipped with the quotient germ universal cover topology, has the structure of a riemannian lamination canonically isometric to $``$.
###### Proof.
Let $`{}_{}{}^{}\stackrel{~}{x}`$, $`{}_{}{}^{}\stackrel{~}{y}[[\stackrel{~}{}]]`$ be such that $`lim{}_{}{}^{}\stackrel{~}{x}=lim{}_{}{}^{}\stackrel{~}{y}=x`$. Thus each point is represented by sequences in $`\stackrel{~}{L}`$ that project to sequences $`\{x_\alpha \}`$, $`\{y_\alpha \}L`$ having a common limit $`x`$. Let $`𝖿`$ be a frame field along a transversal $`T`$ containing $`x`$ and which we may assume contains $`\{x_\alpha \}`$ and $`\{y_\alpha \}`$. Then if $`{}_{}{}^{}f`$, $`{}_{}{}^{}g`$ are the diophantine approximations associated to $`\{x_\alpha \}`$, $`\{y_\alpha \}`$ we have $`{}_{}{}^{}u={}_{}{}^{}g{}_{}{}^{}f_{}^{1}[[\pi ]]_1(,L,x,𝖿)`$ identifies $`{}_{}{}^{}\stackrel{~}{x}`$ with $`{}_{}{}^{}\stackrel{~}{y}`$. Since this latter groupoid belongs to the mother germ by Proposition 5, it follows that $`[[\pi ]]_1()\backslash [[\stackrel{~}{}]]`$ contains all of the identifications implied by $`{}_{}{}^{}p`$ and so may be identified with $``$ with its quotient topology. ∎
Let $``$, $`^{}`$ be riemannian laminations with dense leaves $`L,L^{}`$. A map $`[[\stackrel{~}{F}]]:[[\stackrel{~}{}]][[\stackrel{~}{}^{}]]`$ is called standard if it is induced by $`F:\stackrel{~}{L}\stackrel{~}{L}^{}`$ (e.g. compare with the definition found in §3). In addition, $`[[\stackrel{~}{F}]]`$ is called $`[[\pi ]]_1()`$-equivariant if there exists a groupoid homomorphism $`[[F]]_{}:[[\pi ]]_1()[[\pi ]]_1(^{})`$ such that
$$[[\stackrel{~}{F}]]\left({}_{}{}^{}u{}_{}{}^{}\stackrel{~}{x}\right)=[[F]]_{}\left({}_{}{}^{}u\right)[[\stackrel{~}{F}]]\left({}_{}{}^{}\stackrel{~}{x}\right)$$
for all $`{}_{}{}^{}u[[\pi ]]_1()`$ and $`{}_{}{}^{}\stackrel{~}{x}\mathrm{𝖣𝗈𝗆}({}_{}{}^{}u)`$.
###### Theorem 6.
Let $`[[\stackrel{~}{F}]]:[[\stackrel{~}{}]][[\stackrel{~}{}^{}]]`$ be a standard, $`[[\pi ]]_1()`$-equivariant map. Then $`[[\stackrel{~}{F}]]`$ covers a unique map $`F:^{}`$.
###### Proof.
By equivariance, the expression
$$F={}_{}{}^{}p_{}^{}[[\stackrel{~}{F}]]{}_{}{}^{}p_{}^{1}$$
yields a well-defined function $`F:^{}`$, continuous because $`{}_{}{}^{}p,{}_{}{}^{}p_{}^{}`$ are open maps and $`[[\stackrel{~}{F}]]`$ is continuous. ∎
In , functoriality of the fundamental germ was demonstrated only with respect to the restricted class of trained lamination maps. The following theorem shows that the mother germ is considerably more flexable. A lamination covering map is a surjective lamination map which is a covering map when restricted to any leaf.
###### Theorem 7.
Let $`F:^{}`$ be a lamination covering map. Then $`F`$ induces an injective homomorphism of mother germs
$$[[F]]_{}:[[\pi ]]_1()[[\pi ]]_1(^{}).$$
###### Proof.
Let $`L`$ be a dense leaf and let $`\stackrel{~}{F}:\stackrel{~}{L}\stackrel{~}{L}^{}`$ be the leaf universal cover lift. Then by Theorem 4, $`\stackrel{~}{F}`$ induces a standard map
$$[[\stackrel{~}{F}]]:[[\stackrel{~}{}]][[\stackrel{~}{^{}}]].$$
We note that since $`\stackrel{~}{F}`$ is injective, $`[[\stackrel{~}{F}]]`$ is a homeomorphism onto its image lamination. Let $`{}_{}{}^{}u[[\pi ]]_1()`$. Then the map
$$[[F]]_{}\left({}_{}{}^{}u\right):=[[\stackrel{~}{F}]]{}_{}{}^{}u[[\stackrel{~}{F}]]^1,$$
defined on $`[[\stackrel{~}{F}]]\left(\mathrm{𝖣𝗈𝗆}({}_{}{}^{}u)\right)`$, is deck for the germ universal covering $`{}_{}{}^{}p_{}^{}`$. Indeed
$`{}_{}{}^{}p_{}^{}\left([[\stackrel{~}{F}]]{}_{}{}^{}u[[\stackrel{~}{F}]]^1\right)`$ $`=`$ $`F({}_{}{}^{}p{}_{}{}^{}u)[[\stackrel{~}{F}]]^1`$
$`=`$ $`F{}_{}{}^{}p[[\stackrel{~}{F}]]^1`$
$`=`$ $`{}_{}{}^{}p_{}^{}.`$
Thus the map $`[[F]]_{}`$ is an injective groupoid homomorphism, and we are done. ∎
We have the following
###### Corollary 2.
The mother germ $`[[\pi ]]_1()`$ is independent of leaf-wise riemannian metric and smooth structure. In particular, $`[[\pi ]]_1()`$ is a topological invariant.
## 7. The Antenna Lamination
In this section, we will calculate the fundamental germ of the antenna Riemann surface lamination of Kenyon and Ghys : it is distinguished by the unusual property of having dense leaves of both planar and hyperbolic conformal type.
We begin by constructing a graphical model of a dense leaf of the antenna lamination. Let $`T_1`$ be the cross with vertices $`V_1=\{(0,0),(\pm 1,0),(0,\pm 2)\}`$ and edges consisting of the line segments connecting $`(0,0)`$ to each of the other four vertices. Suppose that we have constructed $`T_n`$ meeting the $`x`$-axis in the interval $`[2^n+1,2^n1]\times \{0\}`$ and meeting the $`y`$-axis in the interval $`\{0\}\times [2^n,2^n]`$. Translate $`T_n`$ vertically so that the origin is taken to $`(0,2^n)`$ and consider the images of this translate by rotations of the plane – about the origin – of angles $`0,\pm \pi /2,\pi `$. The union of these images forms a tree; $`T_{n+1}`$ is then obtained by replacing the extremal edges $`[2^{n+1}2,2^{n+1}]\times \{0\}`$ and $`[2^{n+1},2^{n+1}+2]\times \{0\}`$ by $`[2^{n+1}2,2^{n+1}1]\times \{0\}`$ and $`[2^{n+1}+1,2^{n+1}+2]\times \{0\}`$. It follows that $`T_1T_2\mathrm{}`$: we then define
$$T_{\mathrm{}}=\underset{}{lim}T_n.$$
See Figure 1.
Given $`n`$, let $`\mathrm{𝗈𝗋𝖽}_2(n)`$ be the 2-adic order: the largest nonnegative integer $`r`$ for which $`2^r`$ divides $`n`$. Then the vertex set of $`T_{\mathrm{}}`$ is
$$V_{\mathrm{}}=\left\{(0,0)\right\}\left\{v=(x,y)\right|\mathrm{𝗈𝗋𝖽}_2(x)\mathrm{𝗈𝗋𝖽}_2(y)\}.$$
We may view $`V_{\mathrm{}}`$ as a groupoid through its action on itself by addition. In order to avoid confusion, we write $`vw`$ to indicate groupoid composition, in order to distinguish it from the element $`v+w`$.
###### Proposition 6.
For all $`v,wV_{\mathrm{}}`$, the composition $`vw`$ is defined if and only if $`v=w`$.
###### Proof.
Let $`v,wV_{\mathrm{}}`$. We show that $`\mathrm{𝖱𝖺𝗇}(w)=\mathrm{𝖣𝗈𝗆}(v)`$ if and only if $`v=w`$. Suppose $`vw`$. Then we may write
$$v+w=(\underset{\alpha =m}{\overset{M}{}}a_\alpha 2^\alpha ,\underset{\alpha =n}{\overset{N}{}}b_\alpha 2^\alpha )$$
where $`m`$, $`n`$ are the first non-zero indices of the $`2`$-adic expansions of the coordinates. If $`vw`$ is defined, then since $`0\mathrm{𝖣𝗈𝗆}(w)`$, we must have $`v+wV_{\mathrm{}}`$. In particular, at least one of $`m`$ or $`n`$ is nonzero. Suppose it is $`n`$; we may assume without loss of generality that $`m<n`$. Write
$$w=(\underset{\alpha =r}{\overset{R}{}}c_\alpha 2^\alpha ,\underset{\alpha =s}{\overset{S}{}}d_\alpha 2^\alpha ).$$
Let $`x=(0,2^m)`$. If $`r<m`$, then $`x\mathrm{𝖣𝗈𝗆}(w)`$ but $`v+w+xV_{\mathrm{}}`$ i.e. $`\mathrm{𝖱𝖺𝗇}(w)\mathrm{𝖣𝗈𝗆}(v)`$. This is also true when $`rm`$ except for two cases. If $`r>m`$ and $`s=m`$, $`w+x`$ is not defined presisely when $`1=d_s=\mathrm{}=d_{r1}`$ and $`d_r=0`$. Here we take $`x^{}=(2^r,2^m)\mathrm{𝖣𝗈𝗆}(w)`$ and note that $`v+w+x^{}V_{\mathrm{}}`$. If $`r=m`$ and $`s>m`$, then $`w+x`$ is not defined. In this case, it follows from the form of $`v+w`$ that if $`v=(v_1,v_2)`$ then $`\mathrm{𝗈𝗋𝖽}_2(v_1)>m`$, so that $`x\mathrm{𝖣𝗈𝗆}(v)`$. On the other hand, $`xw\mathrm{𝖣𝗈𝗆}(w)`$. Thus $`\mathrm{𝖱𝖺𝗇}(w)\mathrm{𝖣𝗈𝗆}(v)`$ here as well. ∎
The lines $`x=\pm y`$ intersect $`T_{\mathrm{}}`$ at the origin only. Each of the four components of $`T_{\mathrm{}}\{(0,0)\}`$ defines an end, one contained in each of the four components of $`^2\{(x,\pm x)|x\}`$. Equipped with the path metric induced from $`^2`$, $`T_{\mathrm{}}`$ has exactly four orientation preserving isometries, corresponding to the rotations about the origin of angles $`0,\pm \pi /2,\pi `$ (since ends must be taken to ends). On the other hand, $`T_{\mathrm{}}`$ has many partially defined isometries. For example, for $`vV_{\mathrm{}}`$, let $`I_v`$ be the map of $``$ defined
$$I_v(x,y)=v+(x,y).$$
Then there is a maximal subtree $`T_{\mathrm{}}^vT_{\mathrm{}}`$ (not necessarily connected) for which $`I_v(T_{\mathrm{}}^v)T_{\mathrm{}}`$. By definition, $`I_v`$ is isometric on its domain of definition. If $`v`$ has coordinates of large 2-adic order, then $`I_v`$ is defined on a large ball about $`0`$ in $`T_{\mathrm{}}`$. More precisely, if $`v=(x,y)`$ and $`\mathrm{𝗈𝗋𝖽}_2(x),\mathrm{𝗈𝗋𝖽}_2(y)n`$ then $`T_nT_{\mathrm{}}^v`$. Although the inverse $`I_v^1=I_v`$ is always defined at $`0`$, the composition $`I_{v_1}I_{v_2}=I_{v_1+v_2}`$ will not be defined at $`0`$ if $`v_1+v_2V_{\mathrm{}}`$.
We now define a riemannian surface modelled on $`T_{\mathrm{}}`$, which will occur as a dense leaf of the antenna lamination. Regarding $`T_{\mathrm{}}^2\times \{0\}^3`$, it is clear that
$$S_{\mathrm{}}=\text{boundary of a tubular neighborhood of }T_{\mathrm{}}$$
is homeomorphic to a sphere with four punctures. We want to fix a particular realization of $`S_{\mathrm{}}`$ so that the partial isometries $`I_v`$ of $`T_{\mathrm{}}`$ will induce partial isometries of $`S_{\mathrm{}}`$. Torward this end, consider the surfaces shown in Figure 2.
We assume that they are equipped with riemannian metrics and boundary parametrizations so that given any pair of such surfaces and a choice of boundary component of each, the glueings are canonical and isometric. Each riemannian surface corresponds to a subgraph of $`T_{\mathrm{}}`$, and we may build $`S_{\mathrm{}}`$ from these riemannian surfaces using $`T_{\mathrm{}}`$ as a template. The metrics on the building blocks will also be chosen so that when $`S_{\mathrm{}}`$ is assembled within $`^3`$ it is invariant not only with respect to $`\pi /2`$-rotations about the $`z`$-axis, but also $`\pi `$-rotations about the $`x`$ and $`y`$-axes. We think of $`T_{\mathrm{}}`$ as a spine floating inside the tubular neighborhood bounded by $`S_{\mathrm{}}`$, and we project in the positive vertical direction a copy of $`T_{\mathrm{}}`$ onto $`S_{\mathrm{}}`$. We denote this copy also by $`T_{\mathrm{}}`$, and use the symbols $`0`$ and $`v`$ to denote the origin and a generic element of its vertex set as well. Having constructed $`S_{\mathrm{}}`$ in this way, it is clear that every $`I_v`$ induces a partial isometry of $`S_{\mathrm{}}`$ whose domain is the subsurface (with boundary) of $`S_{\mathrm{}}`$ modelled on $`T_{\mathrm{}}^v`$. We denote this partial isometry $`I_v`$ as well.
Let $`S_{\mathrm{}}^+`$ be the intersection of $`S_{\mathrm{}}`$ with the half plane $`z0`$. The universal cover $`\stackrel{~}{S}_{\mathrm{}}`$ of $`S_{\mathrm{}}`$ is built up from “tiles” modelled on $`S_{\mathrm{}}^+`$, glued together side by side according to the same pattern one uses to glue ideal quadrilaterals to obtain the hyperbolic plane as the universal cover of the four times punctured sphere. Fix $`\stackrel{~}{0}\stackrel{~}{S}_{\mathrm{}}`$ a base point lying over $`0`$. The deck group of the universal covering map is $`F_3`$, the free group on three generators.
Let $`𝗐_0`$ be the unit vector based $`0`$ which is parallel to the $`x`$-axis and points in the positive direction. Consider the vector field $`𝖶`$ on the vertices of $`T_{\mathrm{}}`$ obtained by parallel translating $`𝗐_0`$ along $`T_{\mathrm{}}`$. Note that the partial isometry of $`S_{\mathrm{}}`$ induced by $`I_v`$, $`vV_{\mathrm{}}`$, takes $`𝗐_0`$ to $`𝗐_v=𝖶(v)`$. This is not true of the rotations by angles $`\pm \pi /2`$ and $`\pi `$.
Let $`𝖣_{\stackrel{~}{0}}\stackrel{~}{S}_{\mathrm{}}`$ be the fundamental domain containing $`\stackrel{~}{0}`$. We lift $`T_{\mathrm{}}`$ to $`𝖣_{\stackrel{~}{0}}`$, then translate it by $`F_3`$ to obtain a (disconnected) graph $`\stackrel{~}{T}_{\mathrm{}}`$ on $`\stackrel{~}{S}_{\mathrm{}}`$. Let $`\stackrel{~}{𝖶}`$ be the vector field defined on the vertices of $`\stackrel{~}{T}_{\mathrm{}}`$ that is the lift of $`𝖶`$. The partial isometry $`I_v`$ lifts to a partial isometry of $`\stackrel{~}{S}_{\mathrm{}}`$ which maps a region of each fundamental domain $`𝖣`$ into $`𝖣`$: we denote this privileged lift by $`I_v`$ as well, and the set of such privileged lifts is denoted $`𝕀`$. In addition, by composing with elements of $`F_3`$, we obtain new partial isometries covering $`I_v:S_{\mathrm{}}S_{\mathrm{}}`$. We denote by $`\stackrel{~}{𝕀}`$ the set of partial isometries of $`\stackrel{~}{S}_{\mathrm{}}`$ obtained in this way. Then $`F_3,𝕀\stackrel{~}{𝕀}`$, and every element of $`F_3`$ commutes with every element of $`𝕀`$.
We are now ready to describe the antenna lamination. Consider first the space $`𝔸`$ of all trees in $`^2`$ whose vertex set contains the origin $`0`$ and lies within $``$. Each tree $`T𝔸`$ is equipped with the path metric induced from $`^2`$. On $`𝔸`$, we consider the metric
$$d(T,T^{})=\mathrm{exp}(n),$$
where $`n`$ is the largest integer such that the ball of radius $`n`$ about $`0`$ in $`T`$ coincides with that about $`0`$ in $`T^{}`$. $`𝔸`$ is a compact metric space, . Two graphs $`T`$ and $`T^{}`$ are termed equivalent if there exists a translation by $`(x,y)`$ such that $`T+(x,y)=T^{}`$.
Now for any tree $`T𝔸`$, the ball of radius 1 about $`0`$ is a tree $`P𝔸`$ all of whose vertices lie in the set $`\{(0,0)\}\{(\pm 1,\pm 1)\}`$. We write $`|P|4`$ for the number of vertices $`v`$ of $`P`$ different from $`0`$. There are 16 possible such $`P`$, and we may decompose $`𝔸`$ into a disjoint union of clopens $`𝔸_P`$, where $`𝔸_P`$ consists of those trees whose unit ball about $`0`$ is $`P`$.
For each $`P`$, we consider in the spirit of Figure 2 a model pointed Riemann surface $`(\mathrm{\Sigma }_P,z_P)`$ homeomorphic to $`𝕊^2(|P|\text{ open disks})`$. We assume as before that each boundary component $`_v\mathrm{\Sigma }_P`$ – labeled by a vertex $`v(0,0)`$ of $`P`$ – has a fixed paramentrization, so that any two may be identified along their boundaries without ambiguity. Define
$$=\left(\left(𝔸_P\times \mathrm{\Sigma }_P\right)\right)/\text{ gluing},$$
where the gluing is performed as follows. Given $`T𝔸_P`$, $`vP`$, the translate $`T+v`$ is in $`𝔸_P^{}`$ for some $`P^{}`$, where $`vP^{}`$. We then glue the boundaries $`_v\mathrm{\Sigma }_P`$ and $`_v\mathrm{\Sigma }_P^{}`$. These gluings are compatible with the trivial lamination structures on the $`𝔸_P\times \mathrm{\Sigma }_P`$ and thus $``$ has the structure of a riemannian surface lamination. Note that there is an embedding $`𝔸`$ induced by $`𝔸_P\times \{z_P\}𝔸_P\times \mathrm{\Sigma }_P`$.
Each leaf $`L`$ corresponds to an equivalence class of graph $`T𝔸`$, embedded in $`L`$ as a spine. Note that $`S_{\mathrm{}}`$ is the leaf corresponding to the class of $`T_{\mathrm{}}`$. Define the antenna lamination $`_{\mathrm{}}`$ to be the closure of $`S_{\mathrm{}}`$ in $``$.
Denote by $`S_nS_{\mathrm{}}`$ the surface (with boundary) modelled on the subgraph $`T_nT_{\mathrm{}}`$. If centered at a vertex $`vV_{\mathrm{}}`$ there is a subgraph isometric to $`T_n`$, it models a subsurface $`S_n(v)`$ containing $`v`$, and the isometry $`I_v`$ maps $`S_n`$ to $`S_n(v)`$.
The closure of $`V_{\mathrm{}}`$ in $`_{\mathrm{}}`$ defines a transversal $`𝖳`$ through $`0S_{\mathrm{}}`$, and the vector field $`𝖶`$ is transversally continuous with respect to the topology of $`𝖳`$. A point $`vV_{\mathrm{}}`$ is transversally close to $`0`$ if and only if its coordinates have large 2-adic order.
We are now ready to calculate the fundamental germ
$$[[\pi ]]_1(_{\mathrm{}},0,𝖿),$$
where $`𝖿`$ is the orthonormal frame field determined by $`𝖶`$.
Define nested sets
$$\stackrel{~}{G}=\{\stackrel{~}{G}_n\}\stackrel{~}{𝕀}\text{ and }G=\{G_n\}𝕀,$$
$`n=0,1,2,\mathrm{}`$, as follows. We say that $`\stackrel{~}{I}\stackrel{~}{𝕀}`$ is $`n`$-close if the domain of $`\stackrel{~}{I}`$ contains the finite tree $`\stackrel{~}{T}_n\stackrel{~}{T}_{\mathrm{}}𝖣_{\stackrel{~}{0}}`$ corresponding to $`T_n`$, and maps it into the fundamental domain containing $`\stackrel{~}{I}(\stackrel{~}{0})`$. Then $`\stackrel{~}{G}_n`$ consists of the set of $`n`$-close maps and $`G_n`$ the $`n`$-close maps in $`𝕀`$. Observe that
$$F_3=\stackrel{~}{G}_n.$$
For $`IG_n`$, $`I^1G_n`$ also. Moreover, if $`I^{}G_m`$ and the composition $`II^{}`$ is defined at $`\stackrel{~}{0}`$, then it belongs to $`G_N`$, for $`N=\mathrm{min}(m,n)`$.
Let
$$[[G]]=\{\{I_{v_\alpha }I_{v_\alpha ^{}}^1\}=\left\{I_{v_\alpha v_\alpha ^{}}\right\}|\left\{I_{v_\alpha }\right\},\left\{I_{v_\alpha ^{}}\right\}𝕀\text{ and converge w.r.t }G\}/,$$
where the relation $``$ is defined by an ultrafilter $`𝔘`$ as in (1). We denote the elements of $`[[G]]`$ by $`I_{}_{}{}^{}v`$ where $`{}_{}{}^{}v{}_{}{}^{}{}_{}{}^{}`$, and regard $`[[G]]`$ as a groupoid with unit space $`[[\stackrel{~}{_{\mathrm{}}}]]`$, in which the domains of elements are taken to be maximal in the sense defined in §4.
###### Proposition 7.
As a set, $`[[G]]`$ may be identified with $`{}_{}{}^{}_{\widehat{2}}^{}{}_{}{}^{}_{\widehat{2}}^{}`$, where
$${}_{}{}^{}_{\widehat{2}}^{}=\{\{n_\alpha \}|\mathrm{𝗈𝗋𝖽}_2(n_\alpha )\mathrm{}\text{ as }\alpha \mathrm{}\}/.$$
Given $`{}_{}{}^{}v,{}_{}{}^{}w{}_{}{}^{}_{\widehat{2}}^{}{}_{}{}^{}_{\widehat{2}}^{}`$, the composition $`I_{}_{}{}^{}vI_{}_{}{}^{}w`$ is defined if and only if $`{}_{}{}^{}v={}_{}{}^{}w`$.
###### Proof.
Any element $`({}_{}{}^{}n_{1}^{},{}_{}{}^{}n_{2}^{}){}_{}{}^{}_{\widehat{2}}^{}{}_{}{}^{}_{\widehat{2}}^{}`$ may be written $`({}_{}{}^{}n_{1}^{},0)(0,{}_{}{}^{}n_{2}^{})`$ which clearly defines an element of $`[[G]]`$. Now consider $`{}_{}{}^{}v,{}_{}{}^{}w{}_{}{}^{}_{\widehat{2}}^{}{}_{}{}^{}_{\widehat{2}}^{}`$, and suppose that $`{}_{}{}^{}v{}_{}{}^{}w`$. The $`2`$-adic order extends to
$$\mathrm{𝗈𝗋𝖽}_2:{}_{}{}^{}_{\widehat{2}}^{}{}_{}{}^{}_{\mathrm{}}^{}=({}_{}{}^{})\{0\}.$$
Let $`{}_{}{}^{}V_{\mathrm{}}^{}{}_{}{}^{}_{\widehat{2}}^{}{}_{}{}^{}_{\widehat{2}}^{}`$ be the subset of pairs $`{}_{}{}^{}u=({}_{}{}^{}u_{1}^{},{}_{}{}^{}u_{2}^{})`$ for which $`\mathrm{𝗈𝗋𝖽}_2({}_{}{}^{}u_{1}^{})\mathrm{𝗈𝗋𝖽}_2({}_{}{}^{}u_{2}^{})`$. We distinguish four cases depending on whether $`{}_{}{}^{}v,{}_{}{}^{}w{}_{}{}^{}V_{\mathrm{}}^{}`$ or not. If $`{}_{}{}^{}v,{}_{}{}^{}w{}_{}{}^{}V_{\mathrm{}}^{}`$ then we may regard each as a class of sequences $`\{v_\alpha \}`$, $`\{w_\alpha \}V_{\mathrm{}}`$. If $`\{m_\alpha \}`$, $`\{n_\alpha \}`$ , $`\{r_\alpha \}`$, $`\{s_\alpha \}`$ are the sequences of indices occurring as in Proposition 6, then there classes $`{}_{}{}^{}m`$, $`{}_{}{}^{}n`$, $`{}_{}{}^{}r`$, $`{}_{}{}^{}s`$ are totally ordered in $`{}_{}{}^{}`$, hence we may assume the representative sequences are. In particular, we may proceed with the same argument as in Proposition 6: the sequences $`\{x_\alpha \}`$, $`\{x_\alpha ^{}\}`$ define elements of $`[[\stackrel{~}{_{\mathrm{}}}]]`$ which may be used to show that the composition $`I_{}_{}{}^{}vI_{}_{}{}^{}w`$ is not defined. Now suppose that $`{}_{}{}^{}v{}_{}{}^{}V_{\mathrm{}}^{}`$ but $`{}_{}{}^{}w{}_{}{}^{}V_{\mathrm{}}^{}`$. This means that both components of $`{}_{}{}^{}v`$ have the same order denoted $`\mathrm{𝗈𝗋𝖽}_2({}_{}{}^{}v)`$. Then there exists $`{}_{}{}^{}x{}_{}{}^{}V_{\mathrm{}}^{}`$ such that $`{}_{}{}^{}w+{}_{}{}^{}x{}_{}{}^{}V_{\mathrm{}}^{}`$ in which the two components of $`{}_{}{}^{}w+{}_{}{}^{}x`$ have order greater than $`\mathrm{𝗈𝗋𝖽}_2({}_{}{}^{}v)`$. Then both components of $`{}_{}{}^{}v+{}_{}{}^{}w+{}_{}{}^{}x`$ have equal order, which implies that $`I_{}_{}{}^{}vI_{}_{}{}^{}w`$ is not defined. The case where $`{}_{}{}^{}v{}_{}{}^{}V_{\mathrm{}}^{}`$ but $`{}_{}{}^{}w{}_{}{}^{}V_{\mathrm{}}^{}`$ is handled similarly. Now suppose $`{}_{}{}^{}v,{}_{}{}^{}w{}_{}{}^{}V_{\mathrm{}}^{}`$. Here there are two subcases. First suppose that the orders of the components of $`{}_{}{}^{}v`$, $`{}_{}{}^{}w`$ are not equal. Denote by $`\mathrm{𝗈𝗋𝖽}_2({}_{}{}^{}v)`$, $`\mathrm{𝗈𝗋𝖽}_2({}_{}{}^{}w)`$ the common order of the components of $`{}_{}{}^{}v`$, $`{}_{}{}^{}w`$. Then if say $`\mathrm{𝗈𝗋𝖽}_2({}_{}{}^{}v)<\mathrm{𝗈𝗋𝖽}_2({}_{}{}^{}w)`$, we define $`{}_{}{}^{}x=(0,{}_{}{}^{}w_{2}^{})`$ where $`{}_{}{}^{}w_{2}^{}`$ is the second component of $`{}_{}{}^{}w`$. Then $`I_{}_{}{}^{}w({}_{}{}^{}x)`$ is defined but $`I_{{}_{}{}^{}v+{}_{}{}^{}w}({}_{}{}^{}x)`$ is not. If $`\mathrm{𝗈𝗋𝖽}_2({}_{}{}^{}v)>\mathrm{𝗈𝗋𝖽}_2({}_{}{}^{}w)`$ then $`I_{}_{}{}^{}v`$ is defined on $`{}_{}{}^{}y=(0,{}_{}{}^{}v_{2}^{})`$ but $`{}_{}{}^{}y{}_{}{}^{}w`$ does not define an element of $`\mathrm{𝖣𝗈𝗆}(I_{}_{}{}^{}w)`$ since it does not converge to the same point in $`[[\stackrel{~}{_{\mathrm{}}}]]`$ as $`{}_{}{}^{}y`$. What remains is the case when $`\mathrm{𝗈𝗋𝖽}_2({}_{}{}^{}v)=\mathrm{𝗈𝗋𝖽}_2({}_{}{}^{}w)`$. If $`{}_{}{}^{}v+{}_{}{}^{}w`$ lies in $`{}_{}{}^{}V_{\mathrm{}}^{}`$ then $`{}_{}{}^{}w\mathrm{𝖣𝗈𝗆}(I_{}_{}{}^{}v)`$ but not in $`\mathrm{𝖱𝖺𝗇}(I_{}_{}{}^{}w)`$. Otherwise, if the norms of the components of $`{}_{}{}^{}v+{}_{}{}^{}w`$ are equal, then $`\mathrm{𝗈𝗋𝖽}_2({}_{}{}^{}v+{}_{}{}^{}w)>\mathrm{𝗈𝗋𝖽}_2({}_{}{}^{}v)=\mathrm{𝗈𝗋𝖽}_2({}_{}{}^{}w)`$. If we let $`{}_{}{}^{}x=(({}_{}{}^{}v+{}_{}{}^{}w)_1,0)`$ then $`{}_{}{}^{}x\mathrm{𝖣𝗈𝗆}(I_{{}_{}{}^{}v+{}_{}{}^{}w})`$ but not in $`\mathrm{𝖣𝗈𝗆}(I_{}_{}{}^{}w)`$ so it cannot be that $`I_{}_{}{}^{}vI_{}_{}{}^{}w=I_{{}_{}{}^{}v+{}_{}{}^{}w}`$. ∎
###### Theorem 8.
As a set
$$[[\pi ]]_1(_{\mathrm{}},0,𝖿)={}_{}{}^{}F_{3}^{}\times [[G]].$$
The composition $`{}_{}{}^{}u{}_{}{}^{}v`$, where $`{}_{}{}^{}v=({}_{}{}^{}x,I_{}_{}{}^{}v)`$ , $`{}_{}{}^{}w=({}_{}{}^{}y,I_{}_{}{}^{}w)`$ is defined if and only if $`{}_{}{}^{}v={}_{}{}^{}w`$.
###### Proof.
Every element $`\stackrel{~}{𝕀}`$ may be written in the form $`I_v\gamma =\gamma I_v`$ for $`vV_{\mathrm{}}`$ and $`\gamma F_3`$. Moreover, if $`\stackrel{~}{I}\stackrel{~}{G}_n`$, then $`IG_n`$. The second statement follows immediately from Proposition 7. ∎
Thus, although $`{}_{}{}^{}F_{3}^{}\times [[G]]`$ is formally a group, $`[[\pi ]]_1(_{\mathrm{}},0,𝖿)`$ is not a group with respect to the groupoid structure defined by its action on $`[[\stackrel{~}{_{\mathrm{}}}]]`$. It has nevertheless a distinguished subgroup isomorphic to $`{}_{}{}^{}F_{3}^{}`$. On the other hand,
###### Theorem 9.
Any two elements $`{}_{}{}^{}v`$ and $`{}_{}{}^{}w`$ of $`[[\pi ]]_1(_{\mathrm{}},0,𝖿)`$ define composable elements of the mother germ $`[[\pi ]]_1(_{\mathrm{}})`$ by restriction of domains.
###### Proof.
Let $`I_{}_{}{}^{}v`$, $`I_{}_{}{}^{}w`$ be the $`[[G]]`$-coordinates of $`{}_{}{}^{}v`$, $`{}_{}{}^{}w`$. If $`{}_{}{}^{}v`$, $`{}_{}{}^{}w`$ and $`{}_{}{}^{}v+{}_{}{}^{}w`$ belong to $`{}_{}{}^{}V_{\mathrm{}}^{}`$ then by restricting $`{}_{}{}^{}w`$ to the leaf $`S_{\mathrm{}}`$ and restricting $`{}_{}{}^{}v`$ to the leaf of $`[[\stackrel{~}{_{\mathrm{}}}]]`$ containing $`{}_{}{}^{}w`$, we obtain composable elements of $`[[\pi ]]_1(_{\mathrm{}})`$. The other cases are handled similarly and are left to the reader. ∎
The lamination $`_{\mathrm{}}`$ has the following property: every leaf $`LS_{\mathrm{}}`$ is conformal to either $``$ or $`^{}`$ = $`\{(0,0)\}`$, . Hence $`_{\mathrm{}}`$ is neither a suspension nor a locally free action of a Lie group. In particular, the antenna lamination is beyond the purview of the definition of $`[[\pi ]]_1`$ found in .
Given any leaf $`L`$ of $`_{\mathrm{}}`$, one can obtain a graphical model $`T`$ of $`L`$ as a limit of a sequence of translations of $`T_{\mathrm{}}`$. One can then repeat the discussion leading up to Theorem 8 for $`L`$. The proof of the following is left to the reader.
###### Theorem 10.
Let $`L_{\mathrm{}}`$ be any leaf, modelled as above on a graph $`T𝔸`$ with vertex set $`V`$. Then for $`vV`$ and $`𝖿`$ constructed using a vector field as above, $`[[\pi ]]_1(_{\mathrm{}},v,𝖿)`$ may be identified with
$${}_{}{}^{}\pi _{1}^{}L\times [[G]],$$
where $`{}_{}{}^{}\pi _{1}^{}L\times \{0\}`$ is a subgroup with respect to the groupoid structure that is $``$ $`1`$ or $`{}_{}{}^{}`$.
###### Corollary 3.
Let $`vV_{\mathrm{}}S_{\mathrm{}}`$ and $`v^{}V^{}L^{}S_{\mathrm{}}`$. Then choosing frame fields as above, the fundamental germs $`[[\pi ]]_1(_{\mathrm{}},v,𝖿)`$ and $`[[\pi ]]_1(_{\mathrm{}},v^{},𝖿^{})`$ are not isomorphic.
###### Proof.
$`[[\pi ]]_1(_{\mathrm{}},v,𝖿)`$ has a nonabelian subgroup whereas $`[[\pi ]]_1(_{\mathrm{}},v^{},𝖿^{})`$ is an abelian groupoid. ∎
## 8. The $`PSL(2,)`$ Anosov Foliation
Let $`\mathrm{\Gamma }PSL(2,)`$ be a discrete group of finite type, possibly with elliptic elements. The quotient $`\mathrm{\Sigma }=^2/\mathrm{\Gamma }`$ is a finite volume hyperbolic surface orbifold. The unit tangent bundle $`T_1\mathrm{\Sigma }`$ is defined to be the quotient $`\mathrm{T}_1^2/\mathrm{\Gamma }`$. Let $`\rho :\mathrm{\Gamma }\mathrm{𝖧𝗈𝗆𝖾𝗈}(𝕊^1)`$ be the representation obtained by extending the action of $`\mathrm{\Gamma }`$ to the boundary of $`^2`$. Then $`\mathrm{T}_1\mathrm{\Sigma }`$ may be identified with the suspension
$$\left(^2\times 𝕊^1\right)/\mathrm{\Gamma }$$
as follows. Given $`(\stackrel{~}{z},t)^2\times 𝕊^1`$, associate $`𝗏_{\stackrel{~}{z}}\mathrm{T}_1^2`$, the vector based at $`\stackrel{~}{z}`$ and tangent to the ray limiting to $`t`$. This association is $`\mathrm{\Gamma }`$-equivariant and descends to the desired homeomorphism. The expression of $`\mathrm{T}_1\mathrm{\Sigma }`$ as a suspension defines a hyperbolic Riemann surface foliation $``$ on $`\mathrm{T}_1\mathrm{\Sigma }`$, which is also a fiber bundle over $`\mathrm{\Sigma }`$ provided that $`\mathrm{\Gamma }`$ has no elliptic points. $``$ is called an Anosov foliation.
In , we worked with a definition of $`[[\pi ]]_1`$ that was available for suspensions such as $``$ formed from fixed point free $`\mathrm{\Gamma }`$. Unfortunately, this hypothesis excluded the most “explicit” of discrete subgroups of $`PSL(2,)`$, the modular group $`\mathrm{\Gamma }=PSL(2,)`$. The definition provided in this paper is clearly available in this case, and we devote the rest of this section to its consideration.
Two elements $`r,s\{\mathrm{}\}𝕊^1`$ are called equivalent if there exists $`APSL(2,)`$ such that $`A(r)=s`$. Every equivalence class $`[r]`$ of extended reals corresponds to a leaf $`L_{[r]}`$ of $``$, and since all $`PSL(2,)`$-orbits in $`𝕊^1`$ are dense, all leaves are dense. If $`L_{[r]}`$ is isomorphic to the punctured hyperbolic disk $`𝔻^{}`$, then $`[r]`$ is quadratic over $``$. Otherwise, $`L_{[r]}`$ is isomorphic to $`^2`$.
Let us consider the leaf $`L=L_{[0]}𝔻^{}`$ covered by $`\times \{0\}`$. Choose $`xL`$ and a transversal $`T`$ through $`x`$ that is a fiber with respect to the projection onto the modular surface $`\mathrm{\Sigma }`$. We assume that the lift $`\stackrel{~}{x}`$ of $`x`$ to $`^2`$ is not an elliptic point for the action of $`PSL(2,)`$. Define $`𝖿`$ to be the lift of a frame on $`\mathrm{\Sigma }`$ based at the projection of $`T`$. As before, we denote by $`\stackrel{~}{𝖿}`$ the lift of $`𝖿`$ to $`\stackrel{~}{T}_0^2`$ and by $`\stackrel{~}{𝖿}_{\stackrel{~}{y}}`$ its value at $`\stackrel{~}{y}\stackrel{~}{T}_0`$. Note that for $`APSL(2,)`$, $`A_{}\stackrel{~}{𝖿}_{\stackrel{~}{x}}=A_{}\stackrel{~}{𝖿}_{\stackrel{~}{y}}`$ if and only if $`A\mathrm{\Gamma }=PSL(2,)`$. A sequence $`\{A_\alpha \}PSL(2,)`$ defines an $`𝖿`$-diophantine approximation $``$ $`(A_\alpha \stackrel{~}{x},0)`$ projects to a sequence in $`T`$ converging to $`x`$ $``$ $`(\stackrel{~}{x},A_\alpha ^1(0))`$ projects to a sequence in $`T`$ converging to $`x`$ $``$ $`A_\alpha ^1(0)0`$ in $`𝕊^1`$. Note that for any sequence $`\{\gamma _{n_\alpha }\}`$ in the deck group of $`^2L`$,
$$\left\{\gamma _n=\left(\begin{array}{cc}1& 0\\ n& 1\end{array}\right)\right|n\},$$
the sequence $`\{\gamma _\alpha A_\alpha \}`$ also defines an $`𝖿`$-diophantine approximation. The fundamental germ $`[[\pi ]]_1(,x,𝖿)`$ is then formed from the associated sequences $`\{A_\alpha B_\alpha ^1\}`$ where $`\{B_\alpha \}`$ is another $`𝖿`$-diophantine approximation.
###### Note 3.
The sequences $`\{b_\alpha /d_\alpha =A_\alpha ^1(0)\}`$ are hyperbolic diophantine approximations, as defined for example in . See for more on this point. In the case at hand, they give bad diophantine approximations of $`0`$ whenever $`b_\alpha \mathrm{}`$, in the sense that it is never true that for some $`c>0`$ and almost all $`\alpha `$,
$$\left|0\frac{b_\alpha }{d_\alpha }\right|<\frac{c}{d_\alpha ^2}.$$
The $`𝖿`$-diophantine approximations are not stable with respect to the operation of inversion. Indeed, let $`r`$ be any real number, $`\{m_\alpha /n_\alpha \}`$ a sequence of rationals (written in lowest terms) converging to $`r`$. Let $`M_\alpha ,N_\alpha `$ be such that $`m_\alpha M_\alpha n_\alpha N_\alpha =1`$. Assume that the indexing is such that $`\alpha +N_\alpha /m_\alpha \mathrm{}`$ as $`\alpha \mathrm{}`$. Then the sequence $`\{X_\alpha \}`$,
(5)
$$X_\alpha =\left(\begin{array}{cc}(\alpha m_\alpha +N_\alpha )& m_\alpha \\ \alpha n_\alpha +M_\alpha & n_\alpha \end{array}\right)PSL(2,),$$
satisfies $`X_\alpha ^1(0)0`$, but $`X_\alpha (0)r`$. Using this fact, we can now show
###### Theorem 11.
As a set,
$$[[\pi ]]_1(,x,𝖿)=PSL(2,{}_{}{}^{}).$$
###### Proof.
Let $`\{A_\alpha \}`$ be any sequence in $`PSL(2,{}_{}{}^{})`$. Then after passing to a subsequence if necessary, we find $`A_\alpha ^1(0)r`$ for some $`r\{\mathrm{}\}`$. Note that $`r`$ is independent of the class of $`\{A_\alpha \}`$ in $`PSL(2,{}_{}{}^{})`$. We may choose $`\{n_\alpha \}`$ so that $`\gamma _{n_\alpha }A_\alpha ^1(0)0`$. Hence
$$\{A_\alpha \}=\{A_\alpha \gamma _{n_\alpha }\}\{\gamma _{n_\alpha }^1\}$$
defines an element of $`[[\pi ]]_1(,x,𝖿)`$. ∎
It is not difficult to see that with respect to its action on the germ universal cover $`[[\stackrel{~}{}]]`$, $`[[\pi ]]_1(,x,𝖿)`$ is not a group. Indeed, the class of the sequence $`\{X_\alpha ^1\}`$, where $`\{X_\alpha \}`$ is the sequence appearing in (5), is not defined on $`\stackrel{~}{L}`$.
## 9. Mapping Class Group of the Algebraic Universal Cover of a Surface
In this section, we use the fundamental germ to prove a Nielsen type theorem for the algebraic universal cover of a closed surface. We begin by recalling a few facts, referring the reader to for details.
Let $`\mathrm{\Sigma }`$ be a closed surface and let $`𝒢=\{G_\alpha \}`$ be the set of all normal finite index subgroups. For each $`G_\alpha `$, there exists a covering $`\sigma _\alpha :\mathrm{\Sigma }_\alpha \mathrm{\Sigma }`$ defined by the condition that $`\pi _1\mathrm{\Sigma }_\alpha `$ maps isomorphically onto $`G_\alpha `$. If $`G_\alpha G_\beta `$, there is a unique covering $`s_{\alpha \beta }:\mathrm{\Sigma }_\alpha \mathrm{\Sigma }_\beta `$ for which $`\sigma _\alpha =\sigma _\beta s_{\alpha \beta }`$. Hence the collection of $`\sigma _\alpha `$ and $`s_{\alpha \beta }`$ forms an inverse system of surfaces by covering maps.
###### Definition 6.
The algebraic universal cover of $`\mathrm{\Sigma }`$ is the inverse limit
$$\widehat{\mathrm{\Sigma }}=\underset{}{lim}\mathrm{\Sigma }_\alpha .$$
If $`\sigma :Z\mathrm{\Sigma }`$ is any finite covering, then $`\sigma `$ lifts to a homeomorphism
$$\widehat{\sigma }:\widehat{Z}\widehat{\mathrm{\Sigma }}.$$
Thus the algebraic universal cover depends only on the type of $`\mathrm{\Sigma }`$ (elliptic, parabolic, hyperbolic). In fact, there are only two non-trivial examples of algebraic universal covers of closed surfaces: that of the torus and that of a surface of hyperbolic type.
The inverse limit
$$\widehat{\pi }_1\mathrm{\Sigma }=\underset{}{lim}\left(\pi _1\mathrm{\Sigma }\right)/G_\alpha $$
is a Cantor group called the profinite completion of $`\pi _1\mathrm{\Sigma }`$. The homomorphism $`i:\pi _1\mathrm{\Sigma }\widehat{\pi }_1\mathrm{\Sigma }`$ induced by the system of projections $`\pi _1\mathrm{\Sigma }\pi _1\mathrm{\Sigma }/G_\alpha `$ has dense image. Define a representation
$`\varsigma :\pi _1\mathrm{\Sigma }\mathrm{𝖧𝗈𝗆𝖾𝗈}(\widehat{\pi }_1\mathrm{\Sigma })`$
$`\varsigma _\gamma (\widehat{g})=\widehat{g}i(\gamma )^1`$
for $`\gamma \pi _1\mathrm{\Sigma }`$ and $`\widehat{g}\widehat{\pi }_1\mathrm{\Sigma }`$. Then we may identify $`\widehat{\mathrm{\Sigma }}`$ with the suspension of $`\varsigma `$:
$$\widehat{\mathrm{\Sigma }}\left(\stackrel{~}{\mathrm{\Sigma }}\times \widehat{\pi }_1\mathrm{\Sigma }\right)/\pi _1\mathrm{\Sigma }.$$
With this identification, we see that $`\widehat{\mathrm{\Sigma }}`$ is a surface lamination with Cantor transversals homeomorphic to $`\widehat{\pi }_1\mathrm{\Sigma }`$, that is, a solenoid. Moreover, it can also be seen from this presentation that every leaf $`L`$ of $`\widehat{\mathrm{\Sigma }}`$ satisfies
$$\pi _1LG_\alpha .$$
However for closed surfaces, $`G_\alpha =1`$, so here, $`L`$ is simply connected. Each leaf $`L`$ is dense and a path-component of $`\widehat{\mathrm{\Sigma }}`$. For every $`\alpha `$, the pre-image of the projection map $`\widehat{\mathrm{\Sigma }}\mathrm{\Sigma }_\alpha `$ is a fiber transversal, homeomorphic to $`\widehat{\pi }_1\mathrm{\Sigma }_\alpha \widehat{G}_\alpha `$.
Now let $`L`$ be a fixed leaf of $`\widehat{\mathrm{\Sigma }}`$.
###### Definition 7.
The leafed mapping class group of $`\widehat{\mathrm{\Sigma }}`$ is
$$\mathrm{𝖬𝖢𝖦}(\widehat{\mathrm{\Sigma }},L)=\mathrm{𝖧𝗈𝗆𝖾𝗈}(\widehat{\mathrm{\Sigma }},L)/,$$
where $``$ is the relation of homotopy of homeomorphisms.
We denote by $`[h]`$ the mapping class associated to a homeomorphism $`h`$.
Let $`G`$ be a group.
###### Definition 8.
The virtual automorphism group of $`G`$ is
$$\mathrm{𝖵𝖺𝗎𝗍}(G)=\{\varphi :HH^{}|\varphi \text{ an isomorphism and }H,H^{}\text{ finite index subgroups of }G\}/,$$
where $`\varphi _1\varphi _2`$ if there exists $`H^{\prime \prime }<G`$ of finite index, contained in $`\mathrm{𝖣𝗈𝗆}(\varphi _1)\mathrm{𝖣𝗈𝗆}(\varphi _2)`$ and such that $`\varphi _1|_{H^{\prime \prime }}=\varphi _2|_{H^{\prime \prime }}`$.
Note that the equivalence relation $``$ is precisely what is needed to make composition of virtual automorphisms well-defined. We point out also that if $`H<G`$ is of finite index, then $`\mathrm{𝖵𝖺𝗎𝗍}(H)\mathrm{𝖵𝖺𝗎𝗍}(G)`$.
###### Theorem 12.
$`\mathrm{𝖬𝖢𝖦}(\widehat{\mathrm{\Sigma }},L)\mathrm{𝖵𝖺𝗎𝗍}(\pi _1\mathrm{\Sigma })`$.
###### Note 4.
We first learned the statement of this theorem in a conversation with D. Sullivan in 1995. The first proof appeared in the thesis of C. Odden .
###### Proof.
Define a homomorphism
$$\mathrm{\Theta }:\mathrm{𝖵𝖺𝗎𝗍}(\pi _1\mathrm{\Sigma })\mathrm{𝖬𝖢𝖦}(\widehat{\mathrm{\Sigma }},L)$$
as follows. Given $`\varphi :G_1G_2`$ an isomorphism of finite index subgroups of $`\pi _1\mathrm{\Sigma }`$, we may find covers $`\sigma _1,\sigma _2:\mathrm{\Sigma }^{}\mathrm{\Sigma }`$ – indexed by $`G_1`$ and $`G_2`$ – so that
$$(\sigma _2)_{}(\sigma _1)_{}^1=\varphi :$$
this follows from the classical Nielsen theorem. Then we define
$$\mathrm{\Theta }(\varphi )=[\widehat{\sigma }_2\widehat{\sigma }_1^1]$$
where for $`i=1,2`$, $`\widehat{\sigma }_i:\widehat{\mathrm{\Sigma }}^{}\widehat{\mathrm{\Sigma }}`$ is the algebraic universal cover lift. If $`G^{}<\mathrm{𝖣𝗈𝗆}(\varphi )`$, then $`\mathrm{\Theta }(\varphi |_G^{})=\mathrm{\Theta }(\varphi )`$, since $`\mathrm{\Theta }(\varphi |_G^{})`$ is defined by the pair $`\sigma _i\sigma `$, $`i=1,2`$, where $`\sigma :\mathrm{\Sigma }^{\prime \prime }\mathrm{\Sigma }^{}`$ is a cover for which $`\sigma _1\sigma `$ is indexed by $`G^{}`$. Thus $`\mathrm{\Theta }`$ is a well-defined homomorphism.
###### Claim 1.
$`\mathrm{\Theta }`$ is onto.
Let $`h:(\widehat{\mathrm{\Sigma }},L)(\widehat{\mathrm{\Sigma }},L)`$ be a homeomorphism. After performing an isotopy, we may arrange that $`h`$ fixes a point $`x`$ and fiber transversal $`T`$ containing $`x`$. Without loss of generality, we may assume that $`T`$ is a fiber transversal over $`\mathrm{\Sigma }`$. Due to the suspension structure, $`T\widehat{\pi }\mathrm{\Sigma }`$: fix this identification so that $`x1`$ and $`LT\pi _1\mathrm{\Sigma }`$. Since $`h(L)=L`$, we obtain a bijection
$$h_{}:\pi _1\mathrm{\Sigma }\pi _1\mathrm{\Sigma }$$
in which $`h(1)=1`$.
Suppose that for each $`G_\alpha <\pi _1\mathrm{\Sigma }`$, $`h_{}|G_\alpha `$ is not homomorphic. This means that for every $`\alpha `$, there exists $`\gamma _\alpha ,\gamma _\alpha ^{}G_\alpha `$ so that
(6)
$$h_{}\left(\gamma _\alpha \gamma _\alpha ^{}\right)h_{}\left(\gamma _\alpha \right)h_{}\left(\gamma _\alpha ^{}\right).$$
Assuming that $`\widehat{\mathrm{\Sigma }}`$ has been equipped with a hyperbolic metric, say lifted from $`\mathrm{\Sigma }`$, then the sequences $`\{\gamma _\alpha \}`$, $`\{\gamma _\alpha ^{}\}`$ define elements of the fundamental germ
$${}_{}{}^{}\gamma ,{}_{}{}^{}\gamma _{}^{}[[\pi ]]_1(\widehat{\mathrm{\Sigma }},x,𝖿)$$
where $`𝖿`$ is a frame field lifted from a frame on $`\mathrm{\Sigma }`$. But this fundamental germ is a subgroup of the mother germ $`[[\pi ]]_1(\widehat{\mathrm{\Sigma }})`$. By Theorem 7, $`h`$ induces a groupoid isomorphism
$$[[h]]_{}:[[\pi ]]_1(\widehat{\mathrm{\Sigma }})[[\pi ]]_1(\widehat{\mathrm{\Sigma }}),$$
and so we must have
$$[[h]]_{}\left({}_{}{}^{}\gamma {}_{}{}^{}\gamma _{}^{}\right)=[[h]]_{}\left({}_{}{}^{}\gamma \right)[[h]]_{}\left({}_{}{}^{}\gamma _{}^{}\right).$$
This contradicts equation (6). Thus $`h_{}`$ defines an isomorphism when restricted to some $`G_\alpha `$, and this isomorphism determines an element $`\varphi \mathrm{𝖵𝖺𝗎𝗍}(\pi _1\mathrm{\Sigma })`$. Note that that $`\varphi `$ does not depend on the isotopy used to ensure $`h(x)=x`$ since the holonomy group of $`\widehat{\mathrm{\Sigma }}`$ at any point is trivial. Choose $`\sigma _i:\mathrm{\Sigma }^{}\mathrm{\Sigma }`$, $`i=1,2`$, so that $`\mathrm{\Theta }(\varphi )=[\widehat{\sigma }_2\widehat{\sigma }_1^1]`$. To simplify notation, we write $`h_0=\widehat{\sigma }_2\widehat{\sigma }_1^1`$.
Recall that since $`\widehat{\mathrm{\Sigma }}`$ is compact with hyperbolic leaves, the germ universal cover of $`\widehat{\mathrm{\Sigma }}`$ is $`{}_{}{}^{}_{}^{2}`$. The homeomorphisms $`h`$ and $`h_0`$ lift to the standard bijections $`{}_{}{}^{}h`$ and $`{}_{}{}^{}h_{0}^{}`$ of $`{}_{}{}^{}_{}^{2}`$ sharing the same equivariance with respect to the action of the mother germ $`[[\pi ]]_1(\widehat{\mathrm{\Sigma }})`$. In particular, they act identically on the set of galaxies of $`{}_{}{}^{}_{}^{2}`$. For this reason, we may choose a germ universal cover topology for $`{}_{}{}^{}_{}^{2}`$ with respect to which both $`{}_{}{}^{}h`$ and $`{}_{}{}^{}h_{0}^{}`$ are homeomorphisms.
Define a homotopy $`{}_{}{}^{}H_{t}^{}`$ from $`{}_{}{}^{}h`$ to $`{}_{}{}^{}h_{0}^{}`$ as follows. For each $`{}_{}{}^{}z{}_{}{}^{}_{}^{2}`$, $`{}_{}{}^{}H_{t}^{}\left({}_{}{}^{}z\right)`$ is the point subdividing the hyperbolic geodesic connecting $`{}_{}{}^{}h\left({}_{}{}^{}z\right)`$ to $`{}_{}{}^{}h_{0}^{}\left({}_{}{}^{}z\right)`$ into the proportion $`t:1t`$. By construction, $`{}_{}{}^{}H_{t}^{}`$ has the same equivariance as $`{}_{}{}^{}h`$ and $`{}_{}{}^{}h_{0}^{}`$ and is in particular continuous. Since its initial and final maps are standard, so is $`{}_{}{}^{}H_{t}^{}`$. By Proposition 6, it descends to a homotopy $`H_t`$ of $`h`$ and $`h_0`$. It follows that $`[h]=[h_0]=\mathrm{\Theta }(\varphi )`$, and $`\mathrm{\Theta }`$ is onto.
###### Claim 2.
$`\mathrm{\Theta }`$ is one-to-one.
If not, then there exists $`\varphi `$ the identity map with $`\mathrm{\Theta }(\varphi )=1`$. But then $`\mathrm{\Theta }(\varphi )`$ would have to induce the identity map on the mother germ; by construction, this can only happen if $`\varphi `$ is trivial. ∎
Let $`\mathrm{𝖬𝗈𝖽}(\widehat{\mathrm{\Sigma }},L)`$ be the Teichmüller modular group of the pair $`(\widehat{\mathrm{\Sigma }},L)`$: the group of homotopy classes of quasiconformal homeomorphisms of $`\widehat{\mathrm{\Sigma }}`$ that preserve $`L`$.
###### Corollary 4.
$`\mathrm{𝖬𝗈𝖽}(\widehat{\mathrm{\Sigma }},L)=\mathrm{𝖬𝖢𝖦}(\widehat{\mathrm{\Sigma }},L)`$.
###### Proof.
This follows from the proof of Theorem 12 and the fact that every finite cover of compact Riemann surfaces is homotopic to a quasiconformal cover. ∎
Theorem 12 can be used to formulate the following conjectural Nielsen-type theorem. Given $`xL`$, the fundamental germ $`[[\pi ]]_1(\widehat{\mathrm{\Sigma }},x)`$ is made up of all sequences $`\{\gamma _\alpha \}`$ converging with respect to the lattice of finite index normal subgroups $`G`$ of $`\pi _1\mathrm{\Sigma }`$, so
$$[[\pi ]]_1(\widehat{\mathrm{\Sigma }},x)\underset{[\pi _1\mathrm{\Sigma }:G]<\mathrm{}}{}{}_{}{}^{}G{}_{}{}^{}\pi _{1}^{}\mathrm{\Sigma }.$$
It follows then that there is a monomorphism $`\mathrm{𝖵𝖺𝗎𝗍}(\pi _1)\mathrm{𝖠𝗎𝗍}([[\pi ]]_1(\widehat{\mathrm{\Sigma }},x))`$, which descends to $`\mathrm{𝖵𝖺𝗎𝗍}(\pi _1)\mathrm{𝖮𝗎𝗍}([[\pi ]]_1(\widehat{\mathrm{\Sigma }},x))`$ upon passage to the quotient. This latter map is also a monomorphism, since no nontrivial virtual automorphism $`\varphi `$ can induce on $`[[\pi ]]_1(\widehat{\mathrm{\Sigma }},x)`$ an inner automorphism. For otherwise, $`\varphi `$ would have to be inner on some subgroup $`H`$, hence all, which is only possible if $`\varphi `$ is trivial. In view of these remarks we
###### Conjecture.
The monomorphism $`\mathrm{𝖬𝖢𝖦}(\widehat{\mathrm{\Sigma }},L)\mathrm{𝖮𝗎𝗍}([[\pi ]]_1(\widehat{\mathrm{\Sigma }},x))`$ is an isomorphism.
We end this section by explaining the importance of Theorem 12 and Corollary 4 in giving a genus independent reformulation the Ehrenpreis conjecture. The classical Ehrenpreis conjecture is:
> Given two closed hyperbolic surfaces $`\mathrm{\Sigma }_1`$ and $`\mathrm{\Sigma }_2`$ and $`ϵ>0`$, there exist finite, locally isometric covering surfaces $`Z_1`$ and $`Z_2`$ of each which are $`(1+ϵ)`$-quasiisometric.
We then have the following equivalent, genus independent version:
> Every orbit of the action of $`\mathrm{𝖬𝗈𝖽}(\widehat{\mathrm{\Sigma }},L)`$ on $`𝒯(\widehat{\mathrm{\Sigma }})`$ is dense.
In other words, the genus independent version says that, although the moduli space
$$𝒯(\widehat{\mathrm{\Sigma }})/\mathrm{𝖬𝗈𝖽}(\widehat{\mathrm{\Sigma }},L)$$
is uncountable, it has the “topology of a point” (i.e. the coarse topology). If affirmed, the Ehrenpreis conjecture would thus provide an explanation for the jump between the existence of moduli (dimension 2) and rigidity (dimension 3 and higher) in hyperbolic geometry. See the articles , for more discussion. |
warning/0506/hep-ph0506079.html | ar5iv | text | # Is Dark Matter Heavy Because of Electroweak Symmetry Breaking? Revisiting Heavy Neutrinos
## 1 Introduction
Stable particles produced thermally in the early universe with weak interactions and mass near $`500`$ GeV can account for the measured density of dark matter. This suggests a possible relationship between dark matter and electroweak physics. Consequently, most efforts to explain the origin of dark matter have focused on the Lightest Super-Partner (LSP) in the Minimal Supersymmetric Standard Model (MSSM) . This possibility is attractive because the presence of weak-scale supersymmetry may resolve the gauge hierarchy problem, predicts gauge coupling unification, and explains the connection between the dark matter scale and weak scale. Yet weak-scale supersymmetry was not discovered at LEPII nor have any conclusive indirect signatures been seen . In light of this, the possibility that nature is seemingly fine-tuned has received renewed interest , and so we reconsider a naive explanation for the existence of weak-scale dark matter that does not rely on supersymmetry: Can it be that the dark matter scale coincides with the electroweak scale because the dark matter particles get their mass entirely from electro-weak symmetry breaking?
In this paper, we show that the answer to this question is yes, that such a framework is consistent with experimental data, and that it does not require elaborate fine-tunings. Moreover, it is surprisingly predictive and suggests resolutions to both the gauge and cosmological constant hierarchy problems within a landscape framework.
In Section 2 we review the evidence for dark matter and briefly consider the experimental constraints on heavy ($`20100`$ GeV) neutrinos produced thermally as dark matter candidates. We conclude that, unlike neutrinos with invariant mass terms, a single neutrino with only a Higgs-Yukawa mass is always under-abundant. We note, however, that suppression of the neutrino coupling to the $`Z`$ boson can remedy the many problems with a single heavy neutrino.
In Section 3 we study the effects of $`Z`$ suppression on the thermal heavy neutrino density. The suppression can be realized by adding to the Standard Model with a Higgs doublet a vector-like “dark sector” consisting of two lepton generations with opposite hypercharge and two neutral singlets (3.1). A chiral symmetry prevents these states from mixing with standard model particles or acquiring invariant mass terms. After electroweak symmetry breaking, the stable $`U(1)_{EM}`$-neutral particle is a mixture of opposite-isospin neutral states, thereby suppressing its coupling to the $`Z`$. We present calculations of the dark matter abundance in this toy model (3.2).
In Section 4, we revisit the model of 3.1 and another two-neutrino model allowing $`Z`$-coupling suppression, focusing on the radiative stability of the near-maximal mixing required to reproduce the measured $`\mathrm{\Omega }_{DM}`$. The second model (4.2) is reminiscent of the particle content of Split Supersymmetry, with the gluino removed and supersymmetric relations badly broken. We also point out in (4.3) that gauge coupling unification can be achieved at $`M_G10^{11}`$ GeV by extending the dark sector to include three of the vector lepton families of Sec. 3.1.
In Section 5, we consider the consistency of the toy models of Sections 3 and 4 with current experimental constraints, and the prospects for detection in the near future in direct (5.1) and indirect (5.2) dark matter searches, as well as in colliders (5.3). We focus on the region of parameter space that is most consistent with the observed dark matter density. We also consider the constraints placed on the model by precision electroweak observables (5.4).
Although our motivation is independent of the weak-Planck hierarchy problem, the models we consider do suggest possible resolutions to this problem within a landscape framework. The crucial link in this argument is Weinberg’s observation that structure formation requires that the cosmological constant be small . By explicitly connecting the dark matter mass to the Higgs vev, the class of models suggested here relates the weak-$`\mathrm{\Lambda }`$ hierarchy to the weak-Planck hierarchy. We also further develop the suggestion in that the gauge hierarchy may be explained using the structure principle because the existence of massive dark matter is tied to the stability of the Higgs vacuum. These two arguments are discussed in Section 6.
## 2 Dark Matter: Evidence and Constraints
The existence of dark matter is by now well established. The most direct evidence comes from comparing the luminous mass of a galaxy or cluster to its total mass, inferred by indirect means. Measurements of galactic rotation curves indicate that the mass of galaxies must be greater and more spread out than the luminous matter they contain, suggesting a dark matter component $`\mathrm{\Omega }_{DM}0.1`$. Various means of estimating the mass of galaxy clusters, such as velocity dispersions and gravitational lensing, imply $`\mathrm{\Omega }_{DM}0.20.3`$ on the cluster scale. A less direct—but more precise—probe of dark matter is the anisotropy of the cosmic microwave background (CMB). Assuming a $`\mathrm{\Lambda }`$CDM model (in which all unknown energy density is assumed to come from either a cosmological constant or non-relativistic dark matter), current measurements indicate a total matter component of $`\mathrm{\Omega }_Mh^2=0.134\pm 0.006`$, with a much smaller baryonic matter component of $`\mathrm{\Omega }_Bh^2=0.023\pm 0.001`$ (which is consistent with the limits on $`\mathrm{\Omega }_B`$ from the abundances of light elements produced in big-bang nucleosynthesis). Along with MACHO searches, this is the principal observational evidence for the non-baryonic nature of at least some of the dark matter. Finally, detailed modeling of our own galaxy suggest a local dark matter density of $`\rho _{DM}0.3`$ GeV/$`cm^3`$. For more careful discussions of the evidence for dark matter and its properties, we refer the reader to the reviews .
In light of the increasingly precise evidence for a sizable dark matter component, many explanations of its nature and origin have been proposed (see for a review of both particle physics-motivated and purely astrophysical candidates). The elementary particle candidates are either (1) thermal relics, which begin in thermal equilibrium in the early universe but decouple from the plasma as the expansion rate of the universe exceeds their annihilation rates or (2) non-thermal relics, produced out of equilibrium (for example, from the decay of another particle that has itself decoupled from the plasma). The models we propose are of the first class. We further specialize to the case of dark matter particles that become non-relativistic before they decouple.
In the hot early universe, the dark matter particles are kept in chemical and kinetic equilibrium with the surrounding particle bath by pair production, annihilation, and scattering. If their interactions with the bath are sufficiently strong, they remain coupled until after they become non-relativistic. As the universe cools below $`Tm_{DM}`$, their density becomes exponentially Boltzmann-suppressed, quenching the annihilations. This “freeze-out” happens quite suddenly, typically at a temperature of $`T_f\frac{m_{DM}}{25}`$. From this point on, the co-moving particle density remains approximately constant. Thus, we can compare the observed dark matter density today to the predicted abundance, which is very roughly
$$\mathrm{\Omega }_{DM}\frac{0.1pb\times c}{\sigma v}.$$
(1)
Assuming the maximum perturbative cross section $`\frac{1}{m_{DM}^2}`$ for the dark matter particle’s annihilations and requiring consistency of Equation (1) with the observed dark matter component places an upper bound on its mass, $`m_{DM}`$ 300 TeV . For a weak-interaction coupling, the mass scale of interest is near a TeV, coincidentally close to the scale of electroweak symmetry breaking. Weakly interacting particles with mass near a TeV (WIMPs) provide especially intriguing dark matter candidates because of their possible connection to electroweak physics.
One simple WIMP candidate is a stable heavy Dirac or Majorana neutrino. Neutrinos with a Dirac- or Majorana-type invariant mass have small enough annihilation rates to account for the observed dark matter abundance if they are heavier than $`500`$ GeV, though this option seems to be ruled out by direct dark matter searches (see Sec. 5.1). The best-motivated and currently well-studied WIMP candidate is the LSP in supersymmetric extensions of the Standard Model. In these models, the dark matter mass and electroweak symmetry breaking scale are close for a reason—because both are determined by the scale of supersymmetry breaking.
Another possible explanation of this proximity of scales is if, whatever may set the EWSB scale, the symmetry breaking itself sets the scale of the dark matter mass, which comes entirely from a Higgs Yukawa coupling. Neutrinos with Standard Model coupling to the $`Z`$ that get mass from electroweak symmetry breaking are well ruled out by experiments. $`Z`$-pole width measurements at LEP have ruled out neutrinos lighter than $`\frac{M_Z}{2}`$, and a mass above $`\frac{M_Z}{2}`$ would lead to over-annihilation. Thus $`\mathrm{\Omega }_{DM}h^2=0.11\pm .011`$ cannot be obtained from a Standard Model-like neutrino of any mass. We revisit the possibility that the dark matter particle is given mass by electroweak symmetry breaking and observe that the problems associated with heavy neutrino dark matter can be resolved if the particle’s coupling to the $`Z`$ is suppressed.
## 3 Dark Matter with Suppressed $`Z`$-Coupling
### 3.1 A Toy Model
We consider a model suggested in consisting of the minimal Standard Model, a lepton generation ($`L_1`$, $`\overline{E}_1`$) with the usual standard model charges, and one ($`L_2`$,$`\overline{E}_2`$) with opposite hypercharge assignments. To this we add two $`SU(3)\times SU(2)_L\times U(1)_Y`$ singlets ($`s_1`$, $`s_2`$) (this particle content is among the smallest anomaly-free extensions of the Standard Model in which all particles can obtain mass via EW symmetry breaking). A chiral symmetry under which the Standard Model is neutral and the two new lepton doublets have charge opposite that of the four new singlets (or a discrete subgroup thereof) assures that all invariant mass terms vanish and mixing with Standard Model leptons is eliminated, making the lightest particle exactly stable. This chiral symmetry plays a role analogous to that played by R-parity in supersymmetric models.
The most general Lagrangian consistent with this symmetry takes the form
$$_{\text{int}}=Y_1h^cL_1\overline{E}_1+Y_2hL_2\overline{E}_2+(hL_1,h^cL_2)K\left(\begin{array}{c}s_1\\ s_2\end{array}\right),$$
(2)
where $`h^c=i\sigma _2h^{}`$, and implicit products of $`SU(2)`$ doublets are taken to mean $`ABA^Ti\sigma _2B`$. Up to re-phasing, the $`2\times 2`$ matrix $`K`$ can be written as
$$K=\left(\begin{array}{cc}\mathrm{cos}\theta & \mathrm{sin}\theta \\ \mathrm{sin}\theta & \mathrm{cos}\theta \end{array}\right)\left(\begin{array}{cc}\kappa & 0\\ 0& \kappa ^{}\end{array}\right)𝒱,$$
(3)
with $`𝒱`$ unitary.
In terms of mass eigenstates, the Higgs- and $`Z`$-interaction Lagrangian for the $`U(1)_{EM}`$-neutral states takes the form
$$(1+\frac{h}{v})(mNs+m^{}N^{}s^{})+\frac{g}{2\mathrm{cos}\theta _W}Z_\mu (\overline{N},\overline{N}^{})\overline{\sigma }^\mu \left(\begin{array}{cc}\mathrm{cos}(2\theta )& \mathrm{sin}(2\theta )\\ \mathrm{sin}(2\theta )& \mathrm{cos}(2\theta )\end{array}\right)\left(\begin{array}{c}s\\ s^{}\end{array}\right),$$
(4)
where
$$\left(\genfrac{}{}{0pt}{}{N}{N^{}}\right)=\left(\genfrac{}{}{0pt}{}{N_1\mathrm{cos}(\theta )N_2\mathrm{sin}(\theta )}{N_1\mathrm{sin}(\theta )+N_2\mathrm{cos}(\theta )}\right),\left(\genfrac{}{}{0pt}{}{s}{s^{}}\right)=𝒱\left(\genfrac{}{}{0pt}{}{s_1}{s_2}\right),$$
(5)
and $`m=\frac{v}{\sqrt{2}}\kappa `$ and $`m^{}=\frac{v}{\sqrt{2}}\kappa ^{}`$ are the neutral particle masses after $`h`$ takes the vev $`h=\frac{1}{\sqrt{2}}\left({\displaystyle \genfrac{}{}{0pt}{}{0}{v}}\right)`$. The strength of the neutral fermions’ coupling to the $`Z`$ is suppressed by $`ϵ\mathrm{cos}(2\theta )`$ compared to that of a Standard Model neutrino. As we shall see, in the $`\theta \frac{\pi }{4}`$ regime, the particle is a viable dark matter candidate. Henceforth, we refer to the lighter neutral state as a Dirac fermion $`\chi `$, and the heavier as $`\chi ^{}`$.
### 3.2 Relic Density
The framework for calculating relic abundances for thermally produced, stable particles is discussed in Appendix B. For a Standard Model Dirac neutrino, the relic density ignoring all co-annihilations is shown in Figure 1 together with the thermally averaged cross section $`\sigma v`$ at $`x=24x_f`$. Reproducing the observed $`\mathrm{\Omega }_{DM}`$ requires $`\sigma v2`$ pb c (double the estimate from 1 because both $`\chi `$ and $`\overline{\chi }`$ contribute). As can be seen, $`Z`$-mediated annihilations into fermion final states dominate below the $`Z`$-pole so that the relic density coincides with the experimentally observed level only at $`5`$ GeV, a mass which has been ruled out by $`Z`$-width measurements at LEP. Even as the $`Z`$-mediated annihilation rate drops with increasing $`m_{DM}`$, Higgs-mediated annihilations into $`ZZ`$ and $`WW`$ pairs are strong enough that the relic density never exceeds $`10^2`$. Thus, a neutrino with Higgs-Yukawa-type mass can only be the dark matter if the $`Z`$-coupling is suppressed.
The vector lepton toy model defined in 3.1 allows this suppression as does another model introduced in Sec. 4.2. The $`Z`$-suppression in both cases results, as discussed above, from mass-mixing between gauge eigenstates with opposite isospin. There is always a corresponding enhancement of the $`Z`$-coupling between the heavy and light eigenstates, but if one is significantly heavier than the other this has little effect on the relic density.
We have calculated the relic density of the lightest dark sector neutrino as a function of its mass $`m_{DM}`$ and $`Z`$-coupling suppression factor $`ϵ=\mathrm{cos}2\theta `$. Our calculations are based on the perturbative freezeout methods discussed in and the details of our calculation can be found in appendix B. We ignore the effects of co-annihilations, which are unimportant because the dark matter state can be made naturally much lighter than the other states that could co-annihilate with it.
Figures 2, 3, and 4 display our results for Higgs masses of $`m_h=115`$, 140, and 160 GeV respectively. In each case, we consider $`Z`$-suppressions $`ϵ=0.05`$, 0.1, 0.15, and 0.2. The highest curves correspond to the lowest $`ϵ`$, for which annihilation rates are lowest and the resulting relic densities highest.
The viable neutrino mass below the $`Z`$ pole is pushed up when the $`Z`$ coupling is suppressed, from $`m_{DM}5`$ GeV for a Standard Model neutrino to $`m_{DM}35`$ GeV with $`ϵ0.1`$. As we discuss in 5.3, the $`Z`$-suppression of $`0.1`$ and phase-space suppression from proximity of $`m_{DM}`$ to $`m_Z/2`$ is enough for a heavy neutrino in this mass range to have avoided detection at LEP, while further suppression is required to evade the null results of many direct WIMP searches.
In the region $`m_{DM}\frac{M_Z}{2}`$, there is a band of masses for which some $`ϵ0.1`$ would produce the observed abundance. This band is bounded by the $`Z`$-pole and by the onset of annihilation into $`W^+W^{}`$ and $`ZZ`$ pairs. This region sits near the $`s`$-channel Higgs pole, so it is also chopped off where the efficient annihilation near the pole suppresses relic density. For a Higgs mass of $`m_h=115`$ GeV, the relic density $`\mathrm{\Omega }h_{DM}^2`$ is consistent with experiment in the range of $`m_{DM}6085`$ GeV. On the $`Z`$-pole and above this region, near $`m_{DM}90`$ GeV, the observed abundance can be reproduced by assuming very near-maximal mixing ($`ϵ0.01`$).
When $`Z`$-suppression is allowed, there is an interesting region of parameter space in which a heavy neutrino can reproduce the $`\mathrm{\Omega }_{DM}`$ inferred from observations. The supression necessary is, however, far greater than one would expect from a generic choice of Yukawa couplings. In the following section, we examine the radiative stability of this minimal mixing and its possible origins.
## 4 Models With Maximally Mixed Neutrinos
In this section, we discuss two minimal extensions of the Standard Model in which maximally mixed neutrinos with no invariant mass accounts for the dark matter. The first is the toy model discussed above. The second has a particle content reminiscent of Split-SUSY but with no supersymmetric relations among parameters.
### 4.1 Vector Lepton Model
As discussed above, the model of Section 3.1 generically mixes neutrino species of opposite isospin, resulting in supperssion of the $`Z`$-coupling. But near-maximal mixing ($`ϵ0.15`$) is necessary to reproduce the observed relic density and to evade exclusion based on current data. We wish to consider the radiative stability of such small $`ϵ`$.
The discrete symmetry under which
$`(h,h^c)`$ $``$ $`(h^c,h)`$ (6)
$`(L_1,L_2)`$ $``$ $`(L_2,L_1)`$ (7)
$`(\overline{E_1},\overline{E_2})`$ $``$ $`(\overline{E_2},\overline{E_1})`$ (8)
$`(s_1,s_2)`$ $``$ $`(s_2,s_1)`$ (9)
implies that the neutral mass eigenstates are maximally mixed Dirac neutrinos and that the charged states are degenerate. It is, of course, broken by Standard Model hypercharge and fermion mass differences. However, the Standard Model breaking does not generate a deviation from maximum mixing at one-loop order. We have not analyzed the two-loop contributions carefully, but estimate roughly that, assuming this symmetry is exactly preserved in the dark sector at $`M_G=10^{16}`$ GeV, the Standard Model breaking generates an $`ϵ0.0010.01`$ at the electroweak scale. If the charged state Yukawa couplings are not exactly degenerate (but $`ϵ=0`$ at $`M_G`$), then the non-degeneracy $`Y_1=Y_2(1+\delta )`$ at $`M_G`$ radiatively generates at the weak scale a $`Z`$ suppression
$$ϵ\frac{3}{16\pi ^2}\mathrm{ln}\left(\frac{M_G}{m_W}\right)\frac{K_1^2+K_2^2}{K_1^2K_2^2}\delta ,$$
where $`K_1`$ and $`K_2`$ are the eigenvalues of $`\kappa `$. For reasonable values of the parameters, we find that $`ϵ\delta `$. These estimates are justified in Appendix A.2.
One can also consider an approximate $`SU(2)`$ symmetry that mixes the lepton doublets and rotates the Higgs. It acts as
$``$ $`=`$ $`(h^c,h)𝒰_L𝒰_R,`$ (10)
$``$ $`=`$ $`(L_1,L_2)𝒰_L𝒰_R.`$ (11)
and can act on the Standard Model quark doublets as $`\overline{Q}=(\overline{u},\overline{d})\overline{Q}𝒰_R`$. The symmetry is thus an extension of the ordinary Standard Model $`SU(2)_R`$ that enforces the custodial SU(2) $`\rho =1`$ relation, though it acts very differently on the vector leptons than on the Standard Model fermions. The only invariant combination of $``$ and $``$ is
$$Tr[(i\sigma _2^{(R)})^T(i\sigma _2^{(L)})]=hL_1+h^cL_2$$
(12)
Unlike the discrete symmetries above, it is broken by $`Y_1`$ and $`Y_2`$ even when they are degenerate, as well as by $`K`$ unless $`K_{1i}=K_{2i}`$ (i.e. $`\theta =\pi /4`$, with $`\chi =\frac{1}{\sqrt{2}}(N_1+N_2)`$ massless). In the limit that this symmetry is exact in the neutral sector, the mass matrix takes the form
$$K=\frac{1}{\sqrt{2}}\left(\begin{array}{cc}\hfill 1& \hfill 1\\ \hfill 1& \hfill 1\end{array}\right)\left(\begin{array}{cc}\kappa (\mu )& 0\\ 0& 0\end{array}\right)\left(\begin{array}{cc}\hfill \mathrm{cos}\alpha (\mu )& \hfill \mathrm{sin}\alpha (\mu )\\ \hfill \mathrm{sin}\alpha (\mu )& \hfill \mathrm{cos}\alpha (\mu )\end{array}\right).$$
(13)
The condition $`m_\chi =0`$ is equivalent to demanding that some linear combination $`s_0\mathrm{sin}\alpha s_1+\mathrm{cos}\alpha s_2`$ of the singlets is non-interacting, and so a mass for $`\chi `$ will never be generated radiatively within the $`SU(3)\times SU(2)_L\times U(1)_Y`$ effective theory. Small explicit breaking is needed and generically results in a $`Z`$-suppression $`ϵ`$ of the same order as $`m_\chi /m_\chi ^{}`$ where $`m_\chi ^{}`$ is the mass of $`\chi ^{}=\frac{1}{\sqrt{2}}(N_1N_2)`$.
### 4.2 An Adjoint/Vector Lepton Model
Another minimal extension of the Standard Model that can naturally implement the electroweak dark matter scenario consists of a vector pair of weak doublets $`L_1`$, $`L_2`$, an $`SU(2)_L`$ adjoint fermion $`T=\left(\begin{array}{cc}\hfill T^0& \hfill T^+\\ \hfill T^{}& \hfill T^0\end{array}\right)`$, and a singlet $`s`$ (this is the particle content of Split-SUSY without a heavy gluino, and with no trace of the supersymmetric relations). As in the vector lepton model, we impose a chiral symmetry to suppress the invariant mass terms and mixing with the Standard Model. The Yukawa Lagrangian is
$$_{int}=Y_1hTL_1+Y_2h^cTL_2+Y_3hL_1+Y_4h^cL_2.$$
(14)
An approximate $`SU(2)_R`$ symmetry acting as
$``$ $`=`$ $`(h^c,h)𝒰_L𝒰_R`$ (15)
$``$ $`=`$ $`(L_1,L_2)𝒰_L𝒰_R,`$ (16)
with $`T`$ and $`s`$ both singlets under $`SU(2)_R`$ requires the Lagrangian to be of the form
$`_{int}`$ $`=`$ $`Y^{}Tr[(i\sigma _2^{(R)})^T(i\sigma _2^{(L)})T]+YTr[(i\sigma _2^{(R)})^T(i\sigma _2^{(L)})]s`$ (17)
$`=`$ $`Y^{}(hTL_1+h^cTL_2)+Ys(hL_1+h^cL_2),`$
requiring $`Y_1=Y_2=Y^{}`$, $`Y_3=Y_4=Y`$. The corresponding low-energy mass matrix is
$``$ $``$ $`{\displaystyle \frac{v}{\sqrt{2}}}(N_1,N_2)\left(\begin{array}{cc}Y& Y^{}\\ Y& Y^{}\end{array}\right)\left(\begin{array}{c}s\\ T^0\end{array}\right){\displaystyle \frac{Y^{}v}{\sqrt{2}}}(T^+E_1+T^{}E_2)`$ (22)
$`=`$ $`(\chi ,\chi ^{})\left(\begin{array}{cc}m& 0\\ 0& m^{}\end{array}\right)\left(\begin{array}{c}s\\ T^0\end{array}\right){\displaystyle \frac{m^{}}{\sqrt{2}}}(T^+E_1+T^{}E_2),`$ (27)
where $`m=Yv`$, $`m^{}=Y^{}v`$, $`\chi =\frac{1}{\sqrt{2}}(N_1+N_2)`$, and $`\chi ^{}=\frac{1}{\sqrt{2}}(N_1N_2)`$. Thus, $`SU(2)_R`$ generates maximally mixed neutrino mass eigenstates $`\chi `$ and $`\chi ^{}`$. Standard Model hypercharge and isodoublet mass splitting break this symmetry, but communicate this breaking at two-loop. If all Yukawa terms are strongly coupled near the scale $`M_G10^{16}`$ GeV, then we obtain a low energy spectrum near $`100140`$ GeV. Because the tree-level charged state masses are smaller than the $`\chi ^{}`$ mass, the dark matter must be the $`\chi s`$ state. To reproduce $`\mathrm{\Omega }_{DM}`$, we expect a $`\chi s`$ mass in the range of $`3545`$ GeV or $`5590`$ GeV. We have not studied the radiative effects in this model, but expect them to be of the same size as those discussed for the discrete symmetries of the previous section.
### 4.3 Gauge Coupling Unification
Both of the models discussed above can easily accommodate gauge unification, though this presents additional difficulties associated with proton decay or large contributions to precision electroweak observables.
If we assume that there is not one vector lepton family but three (one for each Standard Model generation), then gauge coupling unification is achieved quite precisely. Figure 5 shows two-loop gauge coupling unification predictions for $`\alpha _s(M_Z)`$. The predictions are within $`2\sigma `$ of the experimental uncertainty (the uncertainty in our calculation is $`0.004`$), and are in fact slightly closer to the observed $`\alpha _s^{exp}(M_Z)=0.119\pm 0.002`$ than in minimal SU(5) SUSY GUTs that predict $`\alpha _s(M_Z)=0.130\pm 0.004`$ (neglecting thresholds).
However, the predicted unification scale is $`M_G10^{11}`$ GeV, so that avoiding dangerous proton decay rates is difficult in GUT completions of this model. This model also results in a large contribution to the precision electroweak observable $`S`$. With three families we have $`S\frac{1}{\pi }`$, which is nearly 3-$`\sigma `$ from current best fits (see section 5.4).
Adding a heavy SU(3) adjoint fermion (a “gluino”) to the adjoint/vector lepton model recovers the low-energy spectrum of Split-SUSY , and unification is achieved quite precisely with a gluino mass of $`m8`$ TeV and a unification scale of $`M_G10^{16}`$ GeV. As discussed in Sec 5.4, this also results in a large contribution to the $`S`$ parameter.
## 5 Experimental Prospects and Constraints
In this section we discuss the implications of the dark matter scenario of the preceding sections for direct and indirect dark matter searches, collider searches, and precision electroweak observables. Figures 6, 7, and 8 summarize experimental constraints on this model for $`m_H=`$ 115, 140, and 160 GeV. We display the $`1\sigma `$ and $`2\sigma `$ regions consistent with $`\mathrm{\Omega }h^2=0.111\pm 0.006`$ (colored bands), the regions of parameter space excluded by direct $`Z`$-pole observations at LEP (upper left light green and blue curves are 3- and 2-$`\sigma `$ exclusions), and the constraints from direct dark matter searches (processed as described below). The dashed gray contour is the boundary of the DAMA 95% C.L. allowed region , while the three lower curves represent the EDELWEISS (dotted dark green), CRESST (dot-dashed black), and CDMS (dashed red) 90% excluded regions. Anything above the colored bands corresponds to a relic density below the observed total abundance, i.e. a heavy neutrino comprising only a fraction of the dark matter.
The bounds from CDMS imply $`ϵ0.01`$. There are three mass ranges for which the dark matter is not overproduced in this regime: at $`m_{DM}45`$, where even weakly coupled particles annihilate effectively through on-shell $`Z`$’s, at 90-95 GeV, and near the Higgs resonance. In the latter two regions, the dark matter annihilates primarily through the Higgs channel.
### 5.1 Direct Dark Matter Searches
Direct searches for dark matter look for elastic and inelastic collisions between WIMPs in the local dark matter halo and heavy nuclei in detectors. Current searches are most sensitive to spin-independent interactions (which are generally mediated by scalar and vector nucleon currents) and probe the cross sections of these interactions using very heavy nuclei such as Ge, I, W or Xe . Current and future experiments probing spin-independent and spin-dependent WIMP-nucleon scattering include DAMA , CDMS , EDELWEISS , CRESST , PICASSO , SIMPLE , ZEPLIN , XENON , NAIAD , GENIUS and HDMS .
WIMP-nucleon scattering results are usually quoted assuming isospin invariance $`G_pG_n`$, but for $`\chi `$ the coupling to neutrons is $`G_n=\frac{ϵG_F}{2\sqrt{2}}`$, while the proton coupling is greatly suppressed: $`G_p=(14\mathrm{sin}^2(\theta _w))G_n0.076G_n`$ . We naively correct for this by reducing quoted sensitivities by $`\left(\frac{(AZ)0.076Z}{A}\right)^2=`$ $`0.28`$, $`0.30`$, and $`0.32`$ for Ge, I, and W detectors respectively . This reduction is taken into account in the exclusion plots 6 through 8, but the cross-sections in the text are as quoted from the original papers, without the correction for isospin-dependent couplings.
Although to date there has been no conclusive discovery of a WIMP signal, the DAMA experiment (a Na/I detector) has reported a $`6.3\sigma `$ annually modulating signal attributable to a WIMP with mass $`m_{DM}50`$ GeV and WIMP-nucleon cross section $`\sigma 7\times 10^6`$ pb . A mixed heavy neutrino with $`ϵ0.020.04`$ would have a cross-section $`\sigma _{\chi n}10^5`$ pb consistent with the DAMA signal. We see also that there are several masses ($`m_{DM}`$ 40 or 50 GeV) close to the reported DAMA signal mass for which this coupling would reproduce the observed $`\mathrm{\Omega }_{DM}`$.
The DAMA signal seems excluded at the $`99.8\%`$ C.L. by the Ge detectors EDELWEISS and CDMS, which use stronger background discrimination techniques . The bound on WIMP-nucleon cross-section obtained from the CDMS null signal with standard estimates of the halo density and local halo profile is near $`\sigma _{WIMPn}4\times 10^7`$ pb at $`m_{DM}=50`$ Gev (assuming isospin-invariance) . This is well below the cross-section of $`\sigma _{\nu n}2\times 10^2`$ pb for a Standard Model heavy neutrino. Our model is consistent with this CDMS limit if $`ϵ0.01`$ and with the EDELWEISS and CRESST limits if $`ϵ0.02`$. If the one-loop radiative generation of $`ϵ`$ is dominant, we would expect $`0.01ϵ0.1`$ (see Sec. 4.1), and $`\sigma _{\chi n}10^610^4`$ pb. Current experiments are probing this regime, and appear to favor smaller $`ϵ`$ and hence charged states that are very nearly degenerate. Two-loop effects from the Standard Model generate $`0.001ϵ0.01`$ ($`10^8\sigma _{\chi n}few\times 10^6`$ pb), and place a floor on the minimum $`ϵ`$ that can be obtained without resorting to fine-tuning. Current and future direct searches will increase sensitivities to $`\sigma _{WIMPn}10^9`$ pb, and so will detect a signal or rule out the viable parameter space for models of this type.
### 5.2 Indirect Searches
Indirect signals of dark matter annihilations in the Earth, Sun or Galactic halo can also be used to probe dark matter. The decay products of annihilating WIMPs concentrated by gravity can give rise to sizable fluxes of neutrinos, gamma rays and cosmic rays above the background of conventional astrophysical sources. We discuss here the general outlook for indirect searches in the electroweak dark matter scenario. We leave a detailed study of the indirect signatures of heavy maximally mixed neutrino dark matter for a future work.
As a large body travels through the local dark matter halo, collisions between dark matter particles and nuclei in the body can dissipate enough energy that the dark matter particles become gravitationally bound and sink to the center of the body. An equilibrium is eventually expected between the rate of capture and of annihilations enhanced by the greater density of dark matter in the center of the body. Excess high-energy neutrino fluxes are the most reliable indicator of WIMP annihilations in the Sun or Earth, and result in bounds on spin-dependent and spin-independent WIMP-nucleon cross sections. Current and future neutrino telescopes carrying out such observations include Super-Kamiokande , AMANDA , BAIKAL , ANTARES , NESTOR and IceCube . To date, no statistically significant neutrino excesses from the Sun or Earth have been detected .
In addition to neutrinos, gamma rays and high-energy cosmic rays from the galactic center also offer a strong indirect detection possibility . Current and future experiments designed to measure gamma and cosmic ray fluxes include EGRET , HEAT , PAMELA , AMS , and GLAST . An excess in gamma rays above 1 GeV detected by EGRET and in positrons peaking at 8 GeV by HEAT may be products of WIMP annihilations in the galactic halo . Moreover, a possible excess of microwave emission from the galactic center observed in WMAP microwave data is consistent with synchrotron radiation from $`e^+e^{}`$ pairs produced by WIMP annihilation .
The uncertainties in the annihilation rates and predicted flux excess due to uncertainties in the halo profile of the Milky Way make it difficult to reach any conclusion regarding the significance of these signals, but we consider briefly their consistency with the model presented here. Previous authors have studied indirect signals coming from heavy stable 4th generation neutrinos , but most recent work relating the observed gamma and cosmic ray excesses to WIMP annihilations have focused on the MSSM neutralino. The principal tension in explaining the HEAT and microwave signals from neutralino dark matter arises because, being Majorana fermions, their annihilations into light Dirac fermion final states such as $`e^+e^{}`$ are suppressed. As such, large bost factors are required to explain the strength of the HEAT positron signal . Moreover, as most of the electrons and positrons from neutralino annihilations are indirect products, neutralino models tend to produce softer positron spectra, whereas the analysis of WMAP data in may prefer harder spectra.
Our dark matter candidate is somewhat similar to the low-mass Kaluza-Klein state considered in , in that both are Dirac particles. In that case, the HEAT signal could be explained with less need for large boost factors than in the neutralino scenario, and we expect the same to be true for our model. We further expect that, if $`m_{DM}80`$ GeV, our dark matter candidate would produce a harder spectrum of positrons than in the neutralino case, and perhaps be more consistent with the WMAP “haze”. Because our dark matter is Dirac, however, there may also be tension with the null neutrino flux signal from AMANDA. Further work is necessary to compare the predicted gamma and positron spectra in our models to the EGRET, HEAT, and microwave observations.
### 5.3 Collider Detection
The most direct constraint on the dark matter particle itself, the lightest neutrino $`\chi `$, comes from the $`Z`$-width measurement at LEP. The measured partial width to invisible states is $`499\pm 1.5`$ MeV, obtained by subtracting the visible partial widths from the total width and assuming lepton universality. The standard model predicts a partial width of $`167.29\pm 0.07`$ MeV to each species of neutrino, or $`501.81\pm 0.13`$ MeV overall . If $`\chi `$ is lighter than $`m_Z/2`$, then it will contribute a partial width,
$$\mathrm{\Gamma }_{Z\chi \chi }=(167.29\pm 0.07)ϵ^2\sqrt{1\frac{4m_{DM}^2}{m_Z^2}},$$
(28)
suppressed because of the reduced $`Z`$-coupling and a threshold suppression due to $`\chi `$’s finite mass. Regions of exclusion at two- and three-$`\sigma `$ ($`\mathrm{\Gamma }_Z(\chi \overline{\chi })=`$ 0.2 and 1.7 MeV, respectively) are shown on the left side of the exclusion plots 6, 7, and 8. It should be noted that the prediction of $`\mathrm{\Gamma }_{inv}`$ from the Standard Model alone is already $`1.9\sigma `$ above the measured width.
A further constraint on this framework comes from the production of $`\chi \chi ^{}`$ pairs through an off-shell $`Z`$ at LEPII. Though the $`\chi \overline{\chi }Z`$ and $`\chi ^{}\overline{\chi }^{}Z`$ couplings are suppressed by $`ϵ`$, the mixed couplings $`\chi \overline{\chi }^{}Z`$ are proportional to $`\sqrt{1ϵ^2}`$, and hence essentially unsuppressed. The maximum center of mass energy reached at LEP II was $`200`$ GeV, so $`\chi \chi ^{}`$ pairs would be produced if $`m_\chi +m_\chi ^{}200`$ GeV. The $`\chi ^{}`$ would decay in the detector to a stable $`\chi `$ and two jets or two lepton tracks via an intermediate $`Z`$. Thus, we believe that a two-jet or two-lepton track plus missing energy signal would have been seen at LEP II if the above mass condition were met. We are not aware of constraints in the literature on this decay mode, and a more careful analysis of the statistics is required to determine the precise exclusion range. Because it depends on $`m_\chi +m_\chi ^{}`$, this bound is only constraining when combined with the assumption of high-scale perturbativity of the Yukawa couplings, which requires $`m_\chi ^{}150`$ GeV.
Because the dark sector interacts only weakly, and most of the energy in $`pp`$ collisions at the LHC will be in the form of energetic gluons, production at the LHC will be accompanied by considerable backgrounds. The favored masses $`m_{DM}90`$ GeV for the dark matter $`\chi `$ and $`200`$ GeV for the heavier neutral and charged states make these models accessible to the LHC, but it is unclear how long it will take to obtain sufficient statistics to learn anything concrete about a given model .
### 5.4 Precision Electroweak Observables
Any extension of the Standard Model that couples to the $`SU(2)_L\times U(1)_Y`$ sector contributes to the gauge boson self-energies and hence modifies predictions for the oblique correction parameters $`S`$, $`T`$, and $`U`$ . These contributions are particularly large for extensions involving fermions that get mass only from electroweak symmetry breaking, which do not decouple from the Standard Model even when $`m_fm_Z`$.
The T parameter measures the amount of custodial SU(2) breaking that occurs, (i.e. deviations from $`\rho =1`$), and current constraints imply that $`_i\frac{1}{3}\mathrm{\Delta }m_i^2(85)^2`$ GeV<sup>2</sup> at $`95\%`$ CL, where $`\mathrm{\Delta }m^2=m_1^2+m_2^2\frac{4m_1^2m_2^2}{m_1^2m_2^2}\mathrm{ln}\frac{m_1}{m_2}`$ and the sum is over all isodoublets . $`SU(2)_L`$ fermion doublets also generically contribute positively to the $`S`$ parameter. Current fits suggest $`S=0.13\pm 0.10,T=0.17\pm 0.12,U=0.22\pm 0.13`$ .
The contribution of a single-family vector lepton model to $`S`$ is $`S\frac{1}{3\pi }`$ which is sufficiently small to not be strongly at odds with precision electroweak data. Moreover, under the assumption that the vector lepton sector is close to strong coupling near the GUT scale ($`M_G10^{13}`$ GeV in these models), the spectrum is not generically split by enough to contribute largely to $`T`$. Only in cases where there is a light dark matter candidate (such as when the $`SU(2)_R`$ symmetry of section 4.1 is approximately preserved) do we expect a sizable positive contribution to $`T`$. In this case, the experimentally preferred value of $`S`$ is closer to $`0`$, and hence more consistent with the positive contribution to $`S`$. Further analysis is needed to work out the detailed predictions for electroweak observables with a vector lepton family, but this minimal model is not ruled out.
For the adjoint/vector lepton model, however, contributions to $`S`$ are much larger, $`S\frac{1}{\pi }`$. As with the vector model, the contributions to $`T`$ depend on the detailed form of the spectrum but do not have to be large. Based on the $`S`$ parameter alone, this minimal model is ruled out at $`4\sigma `$ (if there is a large contribution to $`T`$, then the exclusion is at $`3\sigma `$). An invariant mass term $`m_T`$ for the adjoint field can be added while maintaining a dark matter mass arising entirely from electroweak symmetry breaking so long as $`m_TM_Z`$, thereby suppressing the adjoint field’s large contribution to $`S`$.
## 6 The Cosmological Constant, Gauge Hierarchy, and the Structure Principle
In this section, we discuss some of the implications of the “structure principle” as defined in . For our purposes, this principle requires that large-scale structure develop in the universe as it cools. In the context of landscape scenarios, this principle is justified by the weak anthropic argument that biological creatures will not develop in a universe that cannot support the development of stars and other large-scale stable structure. We review how this principle was first applied by Weinberg to explain the smallness of the cosmological constant and discuss an application suggested in to a possible resolution to the gauge hierarchy problem in the vector lepton model.
### 6.1 Predicting the cosmological constant.
Weinberg’s argument predicting a cosmological constant that is very small if nonzero begins with the observation that gravitational collapse via Jeans instabilities can occur only after the universe becomes matter dominated. Moreover, linear sub-horizon perturbations to the energy density can only grow when the energy density in the form of a cosmological constant is smaller than the energy density of matter, so the universe must be in this regime when non-linear structures such as galaxies start to form. After matter-radiation equality, initial perturbations $`\frac{\delta \rho }{\rho }`$ scale as the acceleration parameter $`a`$, so non-linear structures begin to form after the universe has expanded by an amount $`(\frac{\delta \rho }{\rho })^1`$ after matter-radiation equality, and we require for structure formation $`\mathrm{\Lambda }\rho _{MR}(\frac{\delta \rho }{\rho })^3`$, where $`\rho _{MR}`$ is the energy density at matter-radiation equality.
As was pointed out in , if the dark matter is dominated by cold relics a standard perturbative freeze-out calculation gives
$$\rho _{MR}\frac{1}{3g_{}}\left(\frac{10^2}{M_{PL}\sigma v}\right)^4,$$
(29)
where $`\sigma v`$ is the thermal average of the annihilation rate and $`g_{}`$ is the number of effective degrees of freedom. Weinberg’s argument then bounds $`\mathrm{\Lambda }^{\frac{1}{4}}`$ as,
$$\mathrm{\Lambda }^{1/4}\frac{(\frac{\delta \rho }{\rho })^{3/4}}{(3g_{})^{1/4}}\frac{10^2}{M_{PL}\sigma v}.$$
(30)
Without any additional reason for $`\mathrm{\Lambda }^{\frac{1}{4}}`$ to be small, this bound should be roughly saturated, so it is a generic prediction of the structure principle.
For heavy (i.e. more massive than $`M_Z`$) weakly interacting CDM particles, the annihilation cross section $`\sigma v\frac{\alpha ^2}{m_{DM}^2}`$. The above argument then implies
$$\mathrm{\Lambda }^{1/4}\frac{(\frac{\delta \rho }{\rho })^{3/4}}{\alpha ^2}\frac{m_{DM}^2}{M_{PL}},$$
(31)
as pointed out in . Empirically, heavy WIMP scenarios lead to $`\mathrm{\Omega }h^2.1(\frac{m_{DM}}{TeV})^2`$, thereby implying that $`m_{DM}`$ is of order a TeV and that $`\mathrm{\Lambda }^{1/4}\frac{v^2}{M_{PL}}`$.
Although weakly interacting TeV-scale dark matter seems to be empirically consistent, there is no *a priori* reason why $`m_{DM}`$ is so much smaller than $`M_{PL}`$ or so close to $`v`$. One simple possibility is that weak scale supersymmetry implies that $`m_{DM}v`$ as would also be the case for any other solution to the hierarchy problem that also contains a dark matter candidate. If we give up the assumption that the solution to the gauge hierarchy problem also explains dark matter, a simpler possibility emerges: If the dark matter particle gets mass only via electroweak symmetry breaking, then $`m_{DM}`$ will naturally be close to $`v`$, and $`\mathrm{\Lambda }^{1/4}\frac{v^2}{M_{PL}}`$ is predicted!
### 6.2 Generating the gauge hierarchy with the structure principle
In addition to providing an explanation for the smallness of the cosmological constant, the structure principle also suggests a possible explanation of the gauge hierarchy in any model with a Standard Model Higgs and a dark matter candidate that gets mass only via electro-weak symmetry breaking, such as the vector lepton model discussed here. As we will show, electroweak symmetry breaking via a negative $`m_H^2`$ much below $`M_{PL}`$ now becomes essential for structure formation. If the Higgs mass parameter $`m_H^2`$ can be scanned in a landscape scenario and electroweak symmetry breaking can only happen in a very small window about its current value, then we will have found a possible explanation for the hierarchy.
To see how this might work, we consider a suggestion first made in . In the SM with a large top Yukawa coupling to the Higgs, the quartic coupling $`\lambda `$ has a negative UV fixed point $`\lambda _{UV}y_t^2`$ and an IR fixed point $`\lambda _{IR}+y_t^2`$. If the Higgs is sufficiently light (i.e. $`\lambda `$ is small), then for energies beyond some threshold $`M_{cross}`$, $`\lambda `$ will be driven negative by RG running and the vacuum becomes meta-stable. If $`|\lambda (\mu )|`$ does not become too large in the UV, the vacuum in the early universe has only a very small amplitude to decay as the universe cools, and the theory is safe from cosmological disaster. Turning this picture around, suppose $`\lambda (M_{})`$ starts negative at the cutoff scale, say $`\lambda (M_{})\lambda _{UV}`$. As the universe cools $`\lambda (\mu )`$ crosses zero at a scale $`M_{cross}`$ exponentially suppressed with respect to the cutoff. Moreover, the low energy universe looks very different depending on $`m_H^2`$.
Case I ($`m_H^20`$): The higgs vacuum centered around zero is exactly stable and electro-weak symmetry is broken only by the QCD quark condensate. Because $`M_W\alpha _2\mathrm{\Lambda }_{QCD}`$, sphaleron transitions operate all the way down to $`\mathrm{\Lambda }_{QCD}`$, biasing zero baryon number and washing out any baryon asymmetry in the universe down to the freeze-out level of $`\frac{n_B}{s}10^{19}`$. (See for a discussion of this effect in more detail). The universe remains essentially radiation dominated and devoid of baryons, thereby prohibiting structure formation. Thus, the structure principle would rule this range of parameters out.
Case II ($`m_H^2M_{cross}^2`$): The Higgs field is heavier than $`M_{cross}`$, so $`\lambda (m_H)<0`$ and the Higgs vacuum is unstable. Fluctuations trigger $`\varphi M_{PL}`$. When the vacuum decays, the universe becomes dominated by a cosmological constant $`\mathrm{\Lambda }M_{PL}^4`$, so no structure can form in this case either.
Case III ($`M_{cross}m_H^20`$): Because electroweak symmetry is broken below the scale $`M_{cross}`$, all of the matter fields can obtain sizable masses and the universe can become matter-dominated allowing structure to form.
Thus, if the dark matter gets mass from electroweak symmetry breaking, the prediction without any fine-tunings is that the Higgs mass should be close to $`M_{cross}`$. To determine how close $`|m_H^2|`$ should be to $`M_{cross}`$ to avoid fine-tuning, we calculate the 1-loop effective potential and look for the largest values of $`|m_H^2|`$ for which a meta-stable minimum develops in $`V(\varphi )_{1loop}`$ at nonzero $`\varphi `$. Following , we approximate $`\lambda (\mu )`$ near $`M_{cross}`$ using an approximate solution to its RGE. Following appendix C, we see that $`\lambda (\mu )b\mathrm{log}(\frac{\mu }{M_{cross}})`$ where $`b`$ is an RG coefficient and,
$$\frac{V(\varphi )_{1loop}}{M_{cross}^4}ϵ^2(\frac{\varphi }{M_{cross}})\frac{b}{2}\mathrm{log}(\frac{\varphi }{M_{cross}})(\frac{\varphi }{M_{cross}})^4,$$
(32)
where $`ϵ=|m_H|/M_{cross}`$. For the Standard Model, $`b0.076`$ and numerical solutions indicate that a stable secondary minimum develops for $`ϵ0.09`$ at $`\varphi 0.4M_{cross}`$. Typical values for slightly smaller $`ϵ`$ are in the range of $`\varphi 0.2M_{cross}`$. Thus, without fine-tuning to a scale below that suggested by the structure principle, we would expect physical Higgs masses in the range of $`(0.10.2)M_{cross}`$ for the Standard Model. With a more careful analysis of the one-loop effective potential, we expect to find that the threshold mass is slightly higher than this.
In the Standard Model with a higgs mass of $`115`$ GeV, we calculated $`M_{cross}`$ for a top mass in the experimentally allowed range of $`169.2188.5`$ GeV. $`M_{cross}7,40,120,1350`$ TeV for top masses of $`m_t=188.5,178.1,174.3,169.2`$ GeV respectively. If the top were $`m_t188`$ GeV, then a higgs mass of $`115`$ GeV would be entirely consistent with this scenario without fine-tuning. A top mass near its central value of $`m_t=178`$ leaves a small hierarchy to deal with.
Turning to the vector lepton model, the additional Yukawa couplings drive $`\lambda (\mu )`$ negative even faster. Thus, we expect to find $`M_{cross}`$ lower for heavier lepton masses. Setting all heavy lepton masses equal to $`M_V`$, we consider the dependence on $`M_V`$ of $`M_{cross}`$ for several values of Higgs mass in figure 9.
The hierarchy between typical values of $`M_{cross}`$ in the range of 10-50 TeV in the Standard Model is eliminated for a light Higgs and a modestly heavy vector lepton spectrum of $`150`$ GeV, but $`\lambda (\mu )`$ runs dangerously negative in the UV.
In the Standard Model, requiring that $`\lambda (\mu )`$ not run so negative that the vacuum should have decayed during the last $`10^{10}`$ years leads to a bound on the Higgs mass of $`m_H115`$ GeV. For the vector model extension, even with a light spectrum near $`100`$ GeV, $`\lambda (\mu )`$ tends to run too negative in the UV. In our analysis, we calculated the scale at which $`\lambda (\mu )`$ becomes sufficiently negative for the vacuum in our universe to have decayed already (assuming $`m_H115`$ GeV). We discuss the details of this requirement in appendix C. For this rough analysis, we require that $`\lambda (\mu )`$ not be less than $`\lambda _{decay}0.13`$.
Figure 10 displays our results. We found that requiring $`\lambda `$ to not run past the stability bound requires $`m_H155`$ GeV. Thus, for the above explanation of the gauge hierarchy to work, new physics must enter below the scale $`M_{decay}`$ to prevent $`\lambda (\mu )`$ from running too negative. This new physics should make the UV fixed point for $`\lambda `$ less negative. For example, new $`SU(3)`$ fermions at an intermediate scale slows the running of $`\alpha _3`$ and hence helps keep the top Yukawa $`y_t`$ small, in turn making the UV fixed point for $`\lambda `$ less negative. Alternatively, gauging the chiral symmetry that was introduced to forbid mass terms with a strongly coupled U(1) that is broken near $`M_{decay}`$ (which requires adding an extra vector lepton generation to eliminate anomalies) drives the charged and neutral state Yukawa couplings smaller in the UV, increasing the UV fixed point of $`\lambda `$.
Finally, the upper bound derived from the requirement of not hitting any Landau pole up to $`M_G10^{13}`$ GeV requires $`m_H165`$ GeV. The exact limit depends on vector lepton masses. Consequently, the self-consistency region for the Higgs mass in our minimal vector lepton model is $`155m_H165`$ GeV.
## 7 Conclusions
We have considered a simple and well-motivated explanation for the origin of dark matter, namely that it consists of particles that get their mass entirely through electroweak symmetry breaking. This framework connects the scale and existence of massive dark matter quite simply to the electroweak scale. Moreover, this explicit connection affords new possibilities for explaining the smallness of the cosmological constant in relation to the weak scale as well as the gauge hierarchy.
In the simplest models studied here, the dark matter candidate is a mixture of two Dirac neutrinos with opposite isospin, and so has a $`Z`$ coupling suppressed by $`ϵ=\mathrm{cos}(2\theta )`$. Several approximate discrete and continuous symmetries can guarantee the radiative stability of $`ϵ0`$. If $`ϵ0.15`$ is generated radiatively at one-loop order from breaking of the approximate symmetry within the dark matter sector, the observed relic density can be accounted for with a dark matter mass in the range of $`3035`$ GeV or $`5595`$ GeV. This is, however, disfavored by the null results of several direct WIMP searches. The range $`ϵ0.010.001`$ is typical if the approximate symmetry is broken only by the Standard Model and $`ϵ`$ is generated at two-loop order. Then the dark matter mass is predicted to be $`m_{DM}45`$ GeV or $`m_{DM}9095`$ GeV and WIMP-neutron spin-independent cross sections are $`\sigma _{WIMPn}10^710^8`$ pb. Current and future experiments sensitive to $`\sigma _{WIMPn}10^8`$ pb will probe this regime in the next few years and either detect a signal or rule out the viable parameter space of this class of models.
We found that gauge couplings in a three-family vector lepton model unify quite well near $`M_G10^{11}`$ GeV. A second model with the particle content of Split-SUSY unifies precisely at $`M_G`$ with a heavy “gluino” mass of $`8`$ TeV. Future work is needed to build realistic “top-down” GUT models. Further work is also needed to examine consistency of the gamma ray and positron spectra across the $`1100`$ GeV range expected from WIMP annihilations in the galactic halo with HEAT and EGRET results, as well as to consider other indirect search possibilities.
## 8 Acknowledgements
P.C.S. and N.T. would especially like to thank Nima Arkani-Hamed for inspiring this project and for many insightful discussions and comments. Additional thanks to John Huth, Rakhi Mahbubani, Lubos Motl, Lisa Randall, Leonardo Senatore, and Christopher Stubbs for helpful comments during the completion of this work. P.C.S. and N.T. are each supported by NDSEG Fellowships.
## Appendix A Renormalization Group Results For Vector Lepton Model
### A.1 Vector Lepton Model RG Equations
We have used the following $`\beta `$-functions in Sections 4 and 6. They are derived from the general expressions in . Here $`t=\mathrm{ln}\mu `$, where $`\mu `$ is the renormalization scale. $`g_1=\sqrt{\frac{5}{3}}g^{}`$ is the coupling of the $`SU(5)`$-normalized $`U(1)_Y`$. We quote two-loop beta functions for the gauge couplings, one-loop for the Yukawa and quartic couplings.
1. Gauge couplings (2-loop)
$`(4\pi )^2{\displaystyle \frac{dg_i}{dt}}=`$ $`g_i^3b_i+{\displaystyle \frac{g_i^3}{16\pi ^2}}[{\displaystyle \underset{j=1}{\overset{3}{}}}B_{ij}g_j^2\mathrm{Tr}(d_i^u𝐲_{𝐮}^{}{}_{}{}^{}𝐲_𝐮+d_i^d𝐲_{𝐝}^{}{}_{}{}^{}𝐲_𝐝+d_i^e𝐲_{𝐞}^{}{}_{}{}^{}𝐲_𝐞)`$
$`d_i^E(Y_1^2+Y_2^2)d_i^S\mathrm{tr}(K^TK)]`$ (33)
with
$$b=\left(\begin{array}{ccc}\frac{53}{10}& \frac{5}{2}& 7\end{array}\right)B=\left(\begin{array}{ccc}\frac{28}{5}& 3& \frac{44}{5}\\ \frac{27}{20}& 14& 12\\ \frac{11}{10}& \frac{9}{2}& 26\end{array}\right)$$
$$d^u=\left(\begin{array}{ccc}\frac{17}{10}& \frac{3}{2}& 2\end{array}\right)d^d=\left(\begin{array}{ccc}\frac{1}{2}& \frac{3}{2}& 2\end{array}\right)$$
$$d^e=d^E=\left(\begin{array}{ccc}\frac{3}{2}& \frac{1}{2}& 0\end{array}\right)d^S=\left(\begin{array}{ccc}\frac{3}{10}& \frac{1}{2}& 0\end{array}\right)$$
2. Yukawa couplings (1-loop):
$`(4\pi )^2{\displaystyle \frac{d\kappa _{11}}{dt}}`$ $`=`$ $`\kappa _{11}\left(Y_2(S){\displaystyle \frac{9}{20}}g_1^2{\displaystyle \frac{9}{4}}g_2^2+{\displaystyle \frac{3}{2}}\kappa _{11}^2+{\displaystyle \frac{3}{2}}\kappa _{12}^2+3\kappa _{21}^2{\displaystyle \frac{3}{2}}Y_1^2\right)+3\kappa _{12}\kappa _{21}\kappa _{22}`$
$`(4\pi )^2{\displaystyle \frac{d\kappa _{12}}{dt}}`$ $`=`$ $`\kappa _{12}\left(Y_2(S){\displaystyle \frac{9}{20}}g_1^2{\displaystyle \frac{9}{4}}g_2^2+{\displaystyle \frac{3}{2}}\kappa _{11}^2+{\displaystyle \frac{3}{2}}\kappa _{12}^2+3\kappa _{22}^2{\displaystyle \frac{3}{2}}Y_1^2\right)+3\kappa _{11}\kappa _{21}\kappa _{22}`$
$`(4\pi )^2{\displaystyle \frac{d\kappa _{21}}{dt}}`$ $`=`$ $`\kappa _{21}\left(Y_2(S){\displaystyle \frac{9}{20}}g_1^2{\displaystyle \frac{9}{4}}g_2^2+3\kappa _{11}^2+{\displaystyle \frac{3}{2}}\kappa _{21}^2+{\displaystyle \frac{3}{2}}\kappa _{22}^2{\displaystyle \frac{3}{2}}Y_2^2\right)+3\kappa _{12}\kappa _{12}\kappa _{22}`$
$`(4\pi )^2{\displaystyle \frac{d\kappa _{22}}{dt}}`$ $`=`$ $`\kappa _{21}\left(Y_2(S){\displaystyle \frac{9}{20}}g_1^2{\displaystyle \frac{9}{4}}g_2^2+3\kappa _{12}^2+{\displaystyle \frac{3}{2}}\kappa _{21}^2+{\displaystyle \frac{3}{2}}\kappa _{22}^2{\displaystyle \frac{3}{2}}Y_2^2\right)+3\kappa _{12}\kappa _{12}\kappa _{21}`$
$`(4\pi )^2{\displaystyle \frac{dY_1}{dt}}`$ $`=`$ $`Y_1\left(Y_2(S){\displaystyle \frac{9}{4}}g_1^2{\displaystyle \frac{9}{4}}g_2^2{\displaystyle \frac{3}{2}}\kappa _{11}^2{\displaystyle \frac{3}{2}}\kappa _{12}^2+{\displaystyle \frac{3}{2}}Y_1^2\right)`$
$`(4\pi )^2{\displaystyle \frac{dY_2}{dt}}`$ $`=`$ $`Y_2\left(Y_2(S){\displaystyle \frac{9}{4}}g_1^2{\displaystyle \frac{9}{4}}g_2^2{\displaystyle \frac{3}{2}}\kappa _{21}^2{\displaystyle \frac{3}{2}}\kappa _{22}^2+{\displaystyle \frac{3}{2}}Y_2^2\right),`$ (34)
with
$$Y_2(S)=\mathrm{Tr}\left(3𝐲_{𝐮}^{}{}_{}{}^{}𝐲_𝐮+3𝐲_{𝐝}^{}{}_{}{}^{}𝐲_𝐝+𝐲_{𝐞}^{}{}_{}{}^{}𝐲_𝐞\right)+Y_1^2+Y_2^2+\mathrm{tr}(K^TK)$$
where $`\mathrm{Tr}`$ is over the three Standard Model generations, and $`\mathrm{tr}`$ over the two dimensions of $`K`$.
3. Higgs quartic (1-loop)
$`(4\pi )^2{\displaystyle \frac{d\lambda }{dt}}=`$ $`12\lambda ^2\left({\displaystyle \frac{9}{5}}g_1^2+g_2^2\right)\lambda +{\displaystyle \frac{9}{4}}\left({\displaystyle \frac{3}{25}}g_1^4+{\displaystyle \frac{2}{5}}g_1^2g_2^2+g_2^4\right)`$
$`+4Y_2(S)\lambda 4H(S),`$ (35)
with
$$H(S)=\mathrm{Tr}\left(3(𝐲_{𝐮}^{}{}_{}{}^{}𝐲_𝐮)^2+3(𝐲_{𝐝}^{}{}_{}{}^{}𝐲_𝐝)^2+(𝐲_{𝐞}^{}{}_{}{}^{}𝐲_𝐞)^2\right)+Y_1^4+Y_2^4+\mathrm{tr}((K^TK)^2)$$
### A.2 Radiative Generation Of $`ϵ`$
We note first that for small $`\mathrm{\Delta }`$, the matrix
$$\left(\begin{array}{cc}\kappa /\sqrt{2}(1+\mathrm{\Delta })& \kappa ^{}/\sqrt{2}(1+\mathrm{\Delta })\\ \kappa /\sqrt{2}(1\mathrm{\Delta })& \kappa /\sqrt{2}(1\mathrm{\Delta })\end{array}\right)$$
(36)
can be diagonalized by a rotation of $`\theta =\pi /4+\frac{\kappa ^2+\kappa ^2}{\kappa ^2\kappa ^2}\mathrm{\Delta }`$ on the left and $`\frac{2\kappa \kappa ^{}}{\kappa ^2\kappa ^2}\mathrm{\Delta }`$ on the right. The maximally mixed Yukawa matrix is of this form with $`\mathrm{\Delta }=0`$, and all radiative breaking effects we will consider can be parameterized as radiative generation of nonzero $`\mathrm{\Delta }`$. Thus, we can immediately relate a splitting $`\mathrm{\Delta }`$ of the Yukawa matrix elements to the associated $`Z`$-coupling suppression $`ϵ=\mathrm{cos}(2\theta )=2\frac{\kappa ^2+\kappa ^2}{\kappa ^2\kappa ^2}\mathrm{\Delta }`$.
For the estimates that follow, we work in the approximation that, although the eigenvalues $`\kappa `$ and $`\kappa ^{}`$ of $`K`$ also run, their fractional change is small enough that it does not greatly alter the running of $`\mathrm{\Delta }`$.
We consider first the radiative effect of a splitting $`Y_1=Y_2(1+\delta )`$ of the charged state Yukawa couplings at $`M_G`$, assuming maximally mixed $`K`$ at $`M_G`$. From the RGEs above,
$$(4\pi )^2\frac{d\mathrm{\Delta }}{dt}=\frac{3}{4}(Y_2^2Y_1^2)\frac{3}{2}Y_2\delta ,$$
(37)
so an order-of-magnitude estimate of $`ϵ`$ at the weak scale is given by
$$ϵ(M_W)\mathrm{ln}(\frac{M_G}{M_W})\frac{1}{16\pi ^2}\frac{\kappa ^2+\kappa ^2}{\kappa ^2\kappa ^2}3Y_2\delta \delta $$
(38)
assuming $`\kappa ^{}/\kappa 1.52`$.
In the case where $`Y_1=Y_2`$ at $`M_G`$ is exact, $`ϵ`$ is not generated radiatively at one-loop, but there is a two-loop contribution to $`\beta (\mathrm{\Delta })`$ from the Standard Model that goes as $`y_t^2g_2^2`$. Naively, we expect
$$ϵ_{SM}(M_Z)\frac{1}{(4\pi )^2}\mathrm{ln}(\frac{M_G}{M_Z})\frac{\kappa ^2+\kappa ^2}{\kappa ^2\kappa ^2}\alpha _t\alpha _2N,$$
(39)
where N accounts for color factors, the prefactor of the diagram in the RGEs, and the factor of 2 in the running of $`ϵ`$. We estimate a radiatively generated $`ϵ0.001`$, though this could easily change by an order of magnitude if $`N`$ is large or small, or because of the running of the parameters. Because this radiative generation appears to put $`ϵ`$ in the region of parameter space preferred to avoid direct detection, this radiative generation merits more careful study.
## Appendix B Calculating Relic Dark Matter Abundances
In this appendix, we review our calculation of the relic density for the heavy stable state $`\chi `$ of the vector lepton model of section 4.1. Our calculation is based on the discussion in , , and .
Consider the evolution of the number density $`n_1`$ and $`n_2`$ for $`\chi `$ (with mass m)and its antiparticle $`\overline{\chi }`$. Making the standard assumptions of and neglecting the co-annihilation effects discussed in , the evolution of the total number density $`n=n_1+n_2`$ is described by the Boltzmann equation,
$$\dot{n}+3Hn=\frac{1}{2}\sigma v_M(n^2n_{eq}^2),$$
(40)
where $`n_{eq}`$ is the equilibrium number density, the Moller velocity $`v_M`$ is defined so that $`v_Mn_1n_2`$ is Lorentz invariant, and $`\sigma v_M`$ is the thermal average of the annihilation cross section calculated using the numerical integral
$$\sigma v_M=\frac{1}{8m^4TK_2^2(m/T)}_{4m^2}^{\mathrm{}}\sigma (s4m^2)\sqrt{s}K_1(\sqrt{s}/T)𝑑s,$$
(41)
where $`K_i`$ are modified Bessel functions of order $`i`$. It is useful to treat the effect of the expansion of the universe implicitly by using the variable $`Y=\frac{n}{s}`$ where $`s`$ is the total entropy density of the universe. We also use a dimensionless temperature variable $`x=\frac{m}{T}`$, so that 40 becomes
$$\dot{Y}=s\sigma v_M(Y^2Y_{eq}^2).$$
(42)
From this we can obtain,
$$\frac{dY}{dx}=\frac{1}{3H}\frac{ds}{dx}\sigma v_M(Y^2Y_{eq}^2),$$
(43)
where the Hubble parameter is $`H=\sqrt{\frac{8}{3}\pi G\rho }`$ with $`G`$ the gravitational constant and $`\rho `$ the energy density of the universe. Using conventional definitions $`\rho =g_{eff}(T)\frac{\pi ^2}{30}T^4`$, and $`s=h_{eff}(T)\frac{2\pi ^2}{45}T^3`$ , we can finally re-cast the Boltzmann equation into the form
$$\frac{dY}{dx}=(\frac{45G}{\pi })^{1/2}\frac{g_{}^{1/2}m}{x^2}\sigma v_M(Y^2Y_{eq}^2),$$
(44)
where
$$g_{}^{1/2}=\frac{h_{eff}}{g_{eff}^{1/2}}(1+\frac{1}{3}\frac{T}{h_{eff}}\frac{dh_{eff}}{dT}).$$
(45)
In our analysis, $`\sigma v_M`$ was calculated numerically using equation 41 and then used to numerically solve the Boltzmann equation 40. Throughout, we assumed that the heavy charged states all had mass $`m_{heavy}=140`$ GeV and the heavy neutral state $`m_\chi ^{}1.3`$ consistent with a GUT-scale strong coupling assumption and sufficient for coannihilations to be negligible. The annihilation channels included in our analysis are listed in Table 1. Feyncalc was used to simplify the trace and Lorentz algebra involved in calculating the Feynman diagrams. Mathematica version 5.1 was used for all numerical calculations.
## Appendix C Higgs Vacuum Stability
In this appendix we briefly review our calculation of the Higgs vacuum stability bound in the vector lepton model of section 4.1. The physics of vacuum decay in relativistic quantum field theories is discussed in and an analysis applied to the Standard Model can be found in .
With the addition of a single vector lepton family coupled to the Standard Model as in section 3.1, the Higgs quartic running is modified to include yukawa contributions from the heavy charged and neutral states,
$$(4\pi )^2\beta _\lambda =12\lambda ^2+4(3y_t^2+2y_E^2+\kappa ^2)\lambda 4(3y_t^4+2y_E^4+\kappa ^4).$$
(46)
For a small Higgs mass and hence small $`\lambda (m_h)`$, the large Yukawa couplings $`y_t`$ and $`y_e`$ will drive $`\lambda (\mu )`$ negative for large $`\mu `$, thereby making the vacuum unstable. We wish to demand that in the lifetime of the universe and neglecting finite-temperature effects, the probability for the vacuum to decay is less than 1. This requirement furnishes a lower bound on $`\lambda `$. Applying the results of , the decay rate per unit volume for the nucleation of a bubble of true vacuum of radius $`R`$ is $`\frac{\mathrm{\Gamma }}{V}(\frac{1}{R})^4e^{S_E}`$, where $`S_E`$ is the Euclidean action of the ”bounce” solution (we are dropping multiplicative factors of order unity). Using $`S_E=\frac{16\pi ^2}{3\lambda (1/R)}`$, we take the probability for vacuum decay $`p_d`$ to be,
$$p_d\frac{1}{H^4}\frac{dR}{R}\left(\frac{1}{R}\right)^4e^{\frac{16\pi ^2}{3\lambda (1/R)}},$$
(47)
where $`H=2.133h\times 10^{42}=1.514\times 10^{42}`$ GeV ($`h=0.71`$) is the Hubble constant and $`\frac{1}{H^4}`$ accounts for the spatial and temporal extent of the universe. For the class of models under study, $`\lambda `$ runs negative and then tracks its approximate UV fixed point. Thus, in practice the integral (47) is dominated at the scale where $`\lambda `$ is most negative, $`M_{nuc}`$ (this is typically the GUT scale $`M_G10^{11}10^{13}`$ GeV in these models). In this case, we take, $`p_dM_G^4e^{\frac{16\pi ^2}{3\lambda (M_G)}}`$ to get
$$|\lambda (M_G)|\frac{4\pi ^2}{3\mathrm{ln}\frac{M_G}{H}}.$$
(48)
Using $`M_G10^{11}`$ GeV, we find $`\lambda _{decay}0.13`$. |
warning/0506/quant-ph0506215.html | ar5iv | text | # Atomic Quantum Memory for Photonic Qubits via Scattering in Cavity QED
## Abstract
We investigate a scheme of atomic quantum memory to store photonic qubits in cavity QED. This is motivated on the recent observation that the quantum-state swapping between a single-photon pulse and a $`\mathrm{\Lambda }`$-type atom trapped in a cavity is ideally realized via scattering for some specific case in the strong coupling cavity regime \[T. W. Chen, C. K. Law, and P. T. Leung, Phys. Rev. A 69, 063810 (2004)\]. We derive a simple formula for calculating the fidelity of this atom-photon swapping for quantum memory. We further propose a feasible method which implements conditionally the quantum memory operation with the fidelity of almost unity even if the atom-photon coupling is not so strong. This method can also be applied to store a photonic entanglement in spatially separated atomic quantum memories.
Combined systems of atoms and photons have been studied extensively to construct quite promising and efficient quantum networks for information processing and communication qnet . In these quantum networks, quantum-state transfer between photons and atoms (matter) and storage of quantum states are particularly important. Then, numerous methods to implement the quantum-state transfer and quantum memory have been proposed and investigated in various manners qtrqm-1 ; qtrqm-2 ; qtrqm-3 ; qtrqm-4 ; qtrqm-5 ; qtrqm-6 ; qtrqm-7 ; qtrqm-8 ; qtrqm-9 ; qtrqm-10 ; qtrqm-11 . The cavity QED is among the promising schemes to realize such quantum-state operations, which utilizes strong interaction between single atoms and photons inside cavities CQED . Specifically, quantum-state transfer and manipulation are made between a single atom and a single-photon pulse through a scattering in an optical cavity. In this Letter, we investigate a scheme of atomic quantum memory to store photonic qubits in cavity QED. This is motivated on the recent observation that the quantum-state swapping between a single-photon pulse and a $`\mathrm{\Lambda }`$-type atom trapped in a cavity is ideally realized via scattering for some specific case in the strong coupling cavity regime producing the maximal phase shift CLL-2004 .
We consider a one-dimensional cavity bounded by two mirrors, one of which is perfectly reflecting while the other is partially transparent. The electromagnetic field is expanded in terms of the continuous modes with wave number $`k`$, which range over the inside of cavity through the outside free space CLL-2004 . A photonic qubit is encoded in the polarization states $`|k_L`$ and $`|k_R`$ of single-photon pulse as
$`|\varphi _\mathrm{p}`$ $`=`$ $`c_L|\overline{k}_L+c_R|\overline{k}_R,`$ (1)
$`|\overline{k}_{L,R}`$ $`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑kf(k)e^{ikt}|k_{L,R},`$ (2)
where $`f(k)`$ is the normalized spectral amplitude, and $`e^{ikt}`$ represents the asymptotic temporal evolution ($`c=\mathrm{}=1`$ unit). On the other hand, a $`\mathrm{\Lambda }`$-type atom is trapped inside the cavity, which has two degenerate ground states $`|L`$ and $`|R`$ and an excited state $`|e`$. Then, an atomic qubit is encoded in the degenerate ground states as
$`|\psi _\mathrm{a}=a_L|L+a_R|R.`$ (3)
The polarization states $`|k_L`$ and $`|k_R`$ are coupled, respectively, to the transitions $`|L|e`$ and $`|R|e`$ of frequency $`\omega _e`$ in the cavity with the dipole couplings
$`g_{L,R}(k)=\lambda _{L,R}\sqrt{\kappa /\pi }e^{i\theta _{L,R}}/(kk_c+i\kappa ),`$ (4)
where $`k_c`$ is the resonant frequency of the cavity, $`\kappa `$ is the leakage rate of the cavity, $`\lambda _L`$ and $`\lambda _R`$ represent the normalized coupling strengths, and $`\theta _L`$ and $`\theta _R`$ are the phase angles from the dipole transition matrix elements. The atom-photon scattering then takes place through these couplings, and the transformation of the atom-photon states is induced asymptotically as
$`𝒯|Lk_L`$ $`=`$ $`T_{LL}(k)|Lk_L+T_{RL}(k)|Rk_R,`$
$`𝒯|Rk_R`$ $`=`$ $`T_{LR}(k)|Lk_L+T_{RR}(k)|Rk_R,`$
$`𝒯|Lk_R`$ $`=`$ $`|Lk_R,𝒯|Rk_L=|Rk_L,`$ (5)
where the basis states are taken as $`|Lk_L|L|k_L`$, and so on. The scattering matrix elements are calculated explicitly in Ref. CLL-2004 as
$`T_{LL}(k)=e^{i\varphi _s(k)}|\xi _L(k)|^2+|\xi _R(k)|^2,`$
$`T_{RR}(k)=e^{i\varphi _s(k)}|\xi _R(k)|^2+|\xi _L(k)|^2,`$
$`T_{LR}(k)=\xi _L^{}(k)\xi _R(k)(e^{i\varphi _s(k)}1),`$
$`T_{RL}(k)=e^{2i(\theta _L\theta _R)}\xi _L^{}(k)\xi _R(k)(e^{i\varphi _s(k)}1),`$ (6)
where $`\xi _{L,R}(k)g_{L,R}(k)/\sqrt{|g_L(k)|^2+|g_R(k)|^2}`$. Here, the bright state acquires a phase shift $`\varphi _s(k)`$ via scattering. This linear transformation $`𝒯`$ with a complex $`\varphi _s(k)`$ is generally non-unitary due to the loss with a rate $`\gamma `$ induced by the spontaneous emission into the environment.
Swapping for qubit memory
It is observed CLL-2004 that the quantum-state swapping between the atom and photon can be made ideally via scattering in the specific case of the $`\mathrm{\Lambda }`$-type atom with equal but opposite dipole matrix elements, i.e., $`\lambda _L=\lambda _R=\lambda `$ and $`e^{i(\theta _L\theta _R)}=1`$, which provides
$`g_L(k)=g_R(k).`$ (7)
For example, we may take the D1 line of sodium with $`|L=|F=1,m_F=1`$, $`|R=|F=1,m_F=1`$, $`|e=|F=1,m_F=0`$. In fact, with the maximal phase shift $`e^{i\varphi _s(k_c)}=1`$ at the resonance for $`\kappa \gamma /\lambda ^20`$, we have the scattering matrix elements as
$`T_{LR}(k_c)=T_{RL}(k_c)=1,T_{LL}(k_c)=T_{RR}(k_c)=0.`$ (8)
Then, the atom-photon swapping is obtained as
$`|\mathrm{\Phi }_{\mathrm{in}}^{(k)}`$ $`=`$ $`(a_L|L+a_R|R)(c_L|k_L+c_R|k_R),`$
$``$
$`|\mathrm{\Phi }_{\mathrm{swap}}^{(k)}`$ $`=`$ $`(c_R|L+c_L|R)(a_R|k_L+a_L|k_R).`$ (10)
We here note that this swapping is made in a reversible way via scattering. Hence it can be applied to implement an atomic quantum memory for the storage of unknown photonic qubits of polarization. The input photonic qubit is stored (written) via scattering, and it is retrieved (read) by injecting another single-photon pulse.
We now evaluate the fidelity of this atom-photon swapping for the specific case of $`g_L(k)=g_R(k)`$ providing $`T_{LR}(k)=T_{RL}(k)`$ and $`T_{LL}(k)=T_{RR}(k)=1T_{LR}(k)`$. Arbitrary atomic and photonic qubits in Eqs. (1), (3) and (Atomic Quantum Memory for Photonic Qubits via Scattering in Cavity QED) may be taken as the initial state $`|\mathrm{\Phi }_{\mathrm{in}}`$. Then, the density operator of the output state via scattering is given by
$`\rho _{\mathrm{out}}`$ $`=`$ $`|\mathrm{\Phi }_{\mathrm{out}}\mathrm{\Phi }_{\mathrm{out}}|+(1\mathrm{\Phi }_{\mathrm{out}}|\mathrm{\Phi }_{\mathrm{out}})|00|,`$ (11)
$`|\mathrm{\Phi }_{\mathrm{out}}`$ $`=`$ $`𝒯|\mathrm{\Phi }_{\mathrm{in}}={\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑kf(k)e^{ikt}|\mathrm{\Phi }_{\mathrm{out}}^{(k)},`$ (12)
$`|\mathrm{\Phi }_{\mathrm{out}}^{(k)}`$ $`=`$ $`T_{LR}(k)|\mathrm{\Phi }_{\mathrm{swap}}^{(k)}+T_{LL}(k)|\mathrm{\Phi }_{\mathrm{in}}^{(k)}.`$ (13)
Here, the term of $`|00|`$ represents the loss due to the spontaneous emission with $`\mathrm{Tr}\rho _{\mathrm{out}}=1`$. The output photon will eventually be absorbed by matter. Then, by taking the trace over the photon states the fidelity to obtain the desired atomic state $`|\psi _{\mathrm{swap}}=c_R|L+c_L|R`$ is given by
$`F=\left[\psi _{\mathrm{swap}}|\mathrm{Tr}_{(k)}\left[|\mathrm{\Phi }_{\mathrm{out}}^{(k)}\mathrm{\Phi }_{\mathrm{out}}^{(k)}|\right]|\psi _{\mathrm{swap}}\right]_f,`$ (14)
where $`\mathrm{Tr}_{(k)}[\rho ]k_L|\rho |k_L+k_R|\rho |k_R`$, and the average of any function $`G(k)`$ of $`k`$ with the weight $`|f(k)|^2`$ is denoted as
$`[G]_f{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑k|f(k)|^2G(k).`$ (15)
This fidelity is calculated by considering $`T_{LR}(k)+T_{LL}(k)=1`$ as
$`F(D)=F(0)+[1F(0)]D`$ (16)
with $`D=|\psi _{\mathrm{swap}}|\psi _\mathrm{a}|^2=|c_R^{}a_L+c_L^{}a_R|^2`$ ($`0D1`$). Hence, we may take the fidelity of swapping as
$`F_{\mathrm{swap}}=F(0)=\left[|T_{LR}(k)|^2\right]_f`$ (17)
irrespective to the choice of initial state. Here it is desired for optimizing the quantum-state transfer via swapping that the spectral width $`\kappa _\mathrm{p}`$ of the photon pulse with $`f(k)`$ should be made sufficiently smaller than the cavity leakage rate $`\kappa `$, as discussed in Ref. CLL-2004 . Then, we have the quite high fidelity $`F_{\mathrm{swap}}|T_{LR}(k_c)|^21`$ with $`\kappa _\mathrm{p}\kappa `$ in the strong coupling regime $`\kappa \gamma /\lambda ^21`$.
Numerically, by using the formula for the phase shift $`\varphi _s(k)`$ CLL-2004 we obtain $`F_{\mathrm{swap}}(\mathrm{Gaussian})=0.975`$ typically with $`\lambda =5\kappa `$, $`\gamma =0.5\kappa `$ ($`\kappa \gamma /\lambda ^2=0.02`$) and $`\omega _e=k_c`$ for the Gaussian $`|f(k)|^2=\mathrm{exp}[(kk_c)^2/\kappa _\mathrm{p}^2]/(\pi ^{1/2}\kappa _\mathrm{p})`$ with $`\kappa _\mathrm{p}=0.1\kappa `$. We also obtain $`F_{\mathrm{swap}}(\mathrm{Lorentzian})=0.887(0.960)`$ with $`\lambda =5\kappa `$ and $`\gamma =0.5\kappa `$ for the Lorentzian $`|f(k)|^2=(\kappa _\mathrm{p}/\pi )[(kk_c)^2+\kappa _\mathrm{p}^2]^1`$ with $`\kappa _\mathrm{p}=0.1\kappa (0.02\kappa )`$. The atomic detuning does not provide a significant effect on the fidelity for $`|\omega _ek_c|\gamma \kappa `$.
Storage and retrieval with conditional measurements
We next consider the sequence of storage and retrieval of photonic qubit, which may appear somewhat different from simply repeating twice the atom-photon swapping. To be general, we relax the condition (7) for the $`\mathrm{\Lambda }`$-type atom, allowing different $`g_L(k)`$ and $`g_R(k)`$. In this situation, as seen below, even for a not so strong atom-photon coupling the almost faithful quantum memory operation can be achieved conditionally by making some projective measurements. Specifically, the initial state is taken as
$`|\mathrm{\Phi }_{\mathrm{in}}=|R|\varphi _{\mathrm{p1}}|\overline{k}_R^{}=|R(c_L|\overline{k}_L+c_R|\overline{k}_R)|\overline{k}_R^{}.`$ (18)
Here, $`|\varphi _{\mathrm{p1}}`$ is the photonic qubit to be stored and then retrieved. The atomic state is initially prepared to be $`|R`$, and the second photon pulse of $`|\overline{k}_R^{}`$ is injected after a time delay $`\tau `$ ($`\kappa _\mathrm{p}^1\kappa ^1`$) to retrieve the stored qubit. We take for definiteness the same profile $`f(k)`$ for the two photon pulses, though this choice is not essential.
After the scattering of the first photon pulse with the atom, the detection of the polarization “$`L`$” is made on the output photon. This polarization detection is represented by a positive operator valued measure
$`\mathrm{\Pi }(k_L)={\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑k\eta (k)|k_Lk_L|`$ (19)
with the quantum efficiency $`0<\eta (k)1`$. \[The dark count is neglected here since it can actually be made rather small. The terms of more than one photon states may also be discarded effectively in $`\mathrm{\Pi }(k_L)`$ in the present process involving a single atom and photon.\] Then, the resultant state is given by
$`\rho _1`$ $`=`$ $`P(k_L)^1\mathrm{Tr}_{\mathrm{p1}}[\mathrm{\Pi }(k_L)𝒯_1|\mathrm{\Phi }_{\mathrm{in}}\mathrm{\Phi }_{\mathrm{in}}|𝒯_1^{}]`$ (20)
$`=`$ $`P(k_L)^1{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑k\eta (k)|f(k)|^2|\mathrm{\Phi }_1^{(k)}\mathrm{\Phi }_1^{(k)}|,`$
where
$`|\mathrm{\Phi }_1^{(k)}=[T_{LR}(k)c_R|L+c_L|R]|\overline{k}_R^{}`$ (21)
by applying Eq. (5) for $`𝒯_1`$. It is noticed in Eq. (21) that the initial photonic qubit is transferred to the atomic qubit with slight modification by the factor $`T_{LR}(k)`$. The loss term of $`|00|`$ is projected out by the photon detection even with $`\eta (k)<1`$ (and the negligible dark count). The success probability of the photon detection is given by $`P(k_L)=[\eta (k)\mathrm{\Phi }_1^{(k)}|\mathrm{\Phi }_1^{(k)}]_f`$, providing the normalization $`\mathrm{Tr}_{\mathrm{ap2}}\rho _1=1`$.
The retrieval of the photonic qubit is implemented by the scattering of the second photon pulse followed by the conditional detection of the atomic state $`|L`$. The resultant output state is given by
$`\rho _{\mathrm{out}}`$ $`=`$ $`P(L)^1L|𝒯_2\rho _1𝒯_2^{}|L`$ (22)
$`=`$ $`{\displaystyle \frac{|T_{LR}(k_c)|^2}{P(k_L)P(L)}}`$
$`\times {\displaystyle _{\mathrm{}}^{\mathrm{}}}dk\eta (k)|f(k)|^2|\varphi _{\mathrm{out}}^{(k)}\varphi _{\mathrm{out}}^{(k)}|,`$
where
$`|\varphi _{\mathrm{out}}^{(k)}`$ $`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑k^{}f(k^{})e^{ik^{}(t\tau )}`$ (23)
$`\times [r_{LR}(k)c_R|k_R^{}+r_{LR}(k^{})c_L|k_L^{}]`$
with
$`r_{LR}(k)`$ $``$ $`T_{LR}(k)/T_{LR}(k_c).`$ (24)
It is found in Eq. (23) that with $`T_{LR}(k)T_{LR}(k_c)`$, i.e., $`r_{LR}(k)1`$ in the vicinity of resonance $`|kk_c|\kappa _\mathrm{p}\kappa `$ the output state is very closed to the desired photon state as retrieval:
$`|\varphi _{\mathrm{out}}^{(k)}|\varphi _{\mathrm{p2}}=(c_L|\overline{k}_L^{}+c_R|\overline{k}_R^{}).`$ (25)
The net success probability of the storage and retrieval, providing the normalization $`\mathrm{Tr}_{\mathrm{p2}}\rho _{\mathrm{out}}=1`$, is given by
$`P(L)P(k_L)=|T_{LR}(k_c)|^2\left[\eta (k)\varphi _{\mathrm{out}}^{(k)}|\varphi _{\mathrm{out}}^{(k)}\right]_f.`$ (26)
Then, by considering Eqs. (22)–(25) the fidelity for this operation of storage and retrieval is evaluated as
$`F(\mathrm{p1}\mathrm{a}\mathrm{p2})={\displaystyle \frac{\left[\eta (k)|\varphi _{\mathrm{p2}}|\varphi _{\mathrm{out}}^{(k)}|^2\right]_f}{\left[\eta (k)\varphi _{\mathrm{out}}^{(k)}|\varphi _{\mathrm{out}}^{(k)}\right]_f}}.`$ (27)
As seen in Eq. (25), we can obtain the fidelity of almost unity with $`r_{LR}(k)1`$ for $`|kk_c|\kappa _\mathrm{p}\kappa `$, even if the atom-photon coupling is not so strong as $`\lambda _{L,R}\kappa \gamma `$ to give $`|T_{LR}(k_c)|^20.1`$. Specifically, by considering that the quantum efficiency may be constant as $`\eta (k)=\eta `$ in the vicinity of resonance, the fidelity is calculated as $`F(\mathrm{p1}\mathrm{a}\mathrm{p2})=F_{\mathrm{qm}}+[1F_{\mathrm{qm}}]|c_L|^4`$ depending on the initial $`|\varphi _{\mathrm{p1}}`$. Then, the fidelity of this quantum memory operation is given by
$`F_{\mathrm{qm}}=|[r_{LR}]_f|^2/[|r_{LR}|^2]_f1(\kappa _\mathrm{p}\kappa ).`$ (28)
This is a quite remarkable feature of the present conditional scheme for quantum memory, without requiring the specific $`\mathrm{\Lambda }`$-type atom satisfying the condition (7) on the dipole couplings. Although some modification is made on the stored atomic state, as seen in Eq. (21), it is nearly compensated by the read-out process, as seen in Eq. (23), realizing the almost faithful retrieval of the initial photonic qubit. The trade-off of the success probability is instead paid to obtain the high fidelity for the general case.
Numerically, we have estimates of the fidelity and the net success probability as $`F_{\mathrm{qm}}=\text{0.995, 0.994, 0.999}`$, and $`P(L)P(k_L)=(\text{0.975, 0.634, 0.248})\eta `$, for $`\lambda _L=\lambda _R=(\text{5 , 1 , 0.5})\kappa `$, respectively with $`\gamma =0.5\kappa `$, $`\omega _e=k_c`$ and $`\kappa _\mathrm{p}=0.1\kappa `$ for the Gaussian form. Similar estimates are also obtained for the Lorentzian form with $`\kappa _\mathrm{p}=0.01\kappa `$. Therefore, the quite high fidelity $`F_{\mathrm{qm}}`$ is really obtained almost independently of the atom-photon coupling $`\lambda _{L,R}/\kappa 0.5`$.
It should be mentioned that the atomic detection of $`|L`$ can be implemented by injecting the third photon pulse. Specifically, by injecting the photon of $`|\overline{k}_{L}^{}{}_{}{}^{\prime \prime }`$ followed by the polarization detection $`\mathrm{\Pi }(k_R^{\prime \prime })`$ on the output photon, the $`|L`$ component of the initial atomic state is transformed to $`|R`$ via scattering with the success probability $`\eta |T_{RL}(k_c)|^2`$ while the $`|R`$ one is projected out.
In a feasible experiment, a sufficiently weak coherent light of $`|\alpha `$ may be used as an actual single-photon source, though the success probability becomes rather small proportional to $`|\alpha |^2`$. The vacuum contribution is projected out conditionally by the detection of the output photon. The contributions of more than one photon states are small enough for $`|\alpha |^21`$.
Storage of 2-qubit entanglement
We can see that the quantum-state transfer via scattering can also be applied to the storage of 2-qubit entanglement. We prepare two atomic memories and a polarization-entangled pair of photon pulses. Each photon pulse is scattered with the atom inside the respective cavity. In this situation, particularly for the ideal case of $`T_{LR}=T_{RL}=1`$ and $`T_{LL}=T_{RR}=0`$, we obtain the swapping between the generic states of atom pair and photon pair, which may be either entangled or separable, as
$`\left(\begin{array}{c}a_{LL}\\ a_{RR}\\ a_{LR}\\ a_{RL}\end{array}\right)_\mathrm{a}\left(\begin{array}{c}c_{LL}\\ c_{RR}\\ c_{LR}\\ c_{RL}\end{array}\right)_\mathrm{p}\left(\begin{array}{c}c_{RR}\\ c_{LL}\\ c_{RL}\\ c_{LR}\end{array}\right)_\mathrm{a}\left(\begin{array}{c}a_{RR}\\ a_{LL}\\ a_{RL}\\ a_{LR}\end{array}\right)_\mathrm{p},`$ (45)
where the basis states are taken as for the atom pair “a” $`(|LL,|RR,|LR,|RL)`$ and for the photon pair “p” $`(|k_Lk_L,|k_Rk_R,|k_Lk_R,|k_Rk_L)`$.
This sort of 2-qubit transfer can readily be applied to the storage of photonic polarization entanglement as
$`c_{LR}|\overline{k}_L|\overline{k}_R^{}+c_{RL}|\overline{k}_R|\overline{k}_L^{}c_{RL}|LR+c_{LR}|RL.`$ (47)
Specifically, by taking the initial state
$`|\mathrm{\Phi }_{\mathrm{in}}=|RR(c_{LR}|\overline{k}_L|\overline{k}_R^{}+c_{RL}|\overline{k}_R|\overline{k}_L^{}),`$ (48)
we obtain the $`(kk^{})`$-component of the output state via scatterings in the cavities 1 and 2 as
$`|\mathrm{\Phi }_{\mathrm{out}}^{(kk^{})}`$ $`=`$ $`[T_{LR}(k)c_{RL}|LR+T_{LR}(k^{})c_{LR}|RL]|k_Lk_L^{}`$ (49)
$`+|RR|\stackrel{~}{\varphi }_{LR}^{(kk^{})},`$
where $`|\stackrel{~}{\varphi }_{LR}^{(kk^{})}=T_{RR}(k^{})c_{LR}|k_Lk_R^{}+T_{RR}(k)c_{RL}|k_Rk_L^{}`$. Then, the fidelity for the unconditional operation of entanglement transfer is evaluated by tracing over the photon states and the environment denoted by $`|00|`$. For any choice of the initial state it is calculated to be bounded as $`F(\mathrm{p}\mathrm{a})|[T_{LR}]_f|^2=F_{\mathrm{swap}}F_{\mathrm{qm}}`$, which approaches unity for the specific case of $`g_L(k)=\pm g_R(k)`$ with $`e^{i\varphi _s(k_c)}=1`$ in the strong coupling limit. Furthermore, we can make actively the photon detection $`\mathrm{\Pi }(k_L)\mathrm{\Pi }(k_L^{})`$ on the ouptput state in Eq. (49), so that the trade-off of the success probability is made to obtain the high fidelity. Then, the transfer of the photonic entanglement to the atomic memories is implemented almost faithfully even for $`\lambda _{L,R}\kappa \gamma `$, and the fidelity is calculated to be the same as the 1-qubit memory,
$`F_{\text{ent-tr}}=F_{\mathrm{qm}}1(\kappa _\mathrm{p}\kappa ).`$ (50)
The entanglement stored in the pair of atomic quantum memories is alternatively retrieved by injecting single-photon pulses (may be separable) to the atomic memories. In a feasible experiment for this entanglement transfer, a type II down-conversion light can be used as the input polarization-entangled photonic qubit.
In summary, we have investigated a scheme of atomic quantum memory to store photonic qubits in cavity QED. Specifically for a $`\mathrm{\Lambda }`$-type atom with equal but opposite dipole matrix elements, which is trapped inside an optical cavity, the quantum-state swapping between a single-photon pulse and the atom is ideally realized via scattering in the strong coupling cavity regime We have derived a simple formula for calculating the fidelity of this atom-photon swapping for quantum memory. We have further proposed a feasible method to implements conditionally the quantum memory operation with the fidelity of almost unity even for not so strong atom-photon couplings, which is applicable for a general $`\mathrm{\Lambda }`$-type atom with degenerate ground states. This method can also be applied to store an photonic entanglement in spatially separated atomic quantum memories.
The authors would like to thank M. Kitano, A. Kitagawa, and K. Ogure for valuable suggestions and comments. This work has been supported by International Communications Foundation (ICF). |
warning/0506/math0506480.html | ar5iv | text | # Preperiodic points of polynomials over global fields
## 1. Global Fields and Local Fields
In this section we present the necessary fundamentals from the theory of local and global fields. We also set some notational conventions for this paper. Although this material is well known to number theorists, we present it for the convenience of dynamicists. See Section B.1 of or Section 4.4 of for more details on global fields and sets of absolute values; see for expositions concerning the local fields $`_v`$.
### 1.1. Global fields and absolute values
Throughout this paper, $`K`$ will denote a global field. That is, $`K`$ is either a number field (i.e., a finite extension of $``$) or a function field over a finite field (i.e., a finite extension of $`𝔽_p(T)`$ for some prime $`p`$). We will write $`M_K`$ for the set of standard absolute values on $`K`$. That is, $`M_K`$ consists of functions $`||_v:K`$ satisfying $`|x|_v0`$ (with equality if and only if $`x=0`$), $`|xy|_v=|x|_v|y|_v`$, and $`|x+y|_v|x|_v+|y|_v`$, for all $`x,yK`$. (We will frequently abuse notation and write $`vM_K`$ when our meaning is clear.) Moreover, the absolute values in $`M_K`$ are chosen to satisfy a product formula, which is to say that for each $`vM_K`$, there is an integer $`n_v1`$ such that for all $`xK^\times `$,
(1)
$$\underset{vM_K}{}|x|_v^{n_v}=1.$$
(Implicit in the product formula is the fact that for any $`xK^\times `$, we have $`|x|_v=1`$ for all but finitely many $`vM_K`$.)
All but finitely many $`vM_K`$ satisfy the ultrametric triangle inequality
$$|x+y|_v\mathrm{max}\{|x|_v,|y|_v\}.$$
(Note that this means $`|n|_v1`$ for all integers $`n`$.) Such $`v`$ are called non-archimedean absolute values; the finitely many exceptions are called archimedean absolute values. If $`K`$ is a function field, then all absolute values are non-archimedean. If $`K`$ is a number field, then there are archimedean absolute values, each of which, when restricted to $``$, is the familiar absolute value $`||`$, commonly denoted $`||_{\mathrm{}}`$. In fact, if $`K`$ is a number field, then
(2)
$$\underset{\begin{array}{c}vM_K,\\ v\text{ archimedean}\end{array}}{}n_v=[K\mathrm{:}].$$
Meanwhile, the non-archimedean absolute values in $`M_K`$ correspond to prime ideals of the ring of integers of $`K`$. For this reason, we frequently refer to the absolute values $`vM_K`$ as primes of $`K`$, even when $`v`$ is archimedean.
If $`v`$ is non-archimedean, then $`|K^\times |_v`$ is a discrete subset of $``$, and we say that $`v`$ is a discrete valuation on $`K`$. In that case, let $`\epsilon (0,1)`$ be the largest absolute value less than $`1`$ attained in $`|K^\times |_v`$, and choose $`\pi _vK`$ with $`|\pi _v|_v=\epsilon `$. Then $`\pi _v`$ is called a uniformizer of $`K`$ at $`v`$, and we have $`|K^\times |_v=\{\epsilon ^m:m\}`$. Moreover, if $`K`$ is a number field, then $`|\pi _v|_v^{n_v}=p^f`$ for some prime number $`p`$ and some positive integer $`f`$, and $`||_v`$ restricted to $``$ is the usual $`p`$-adic absolute value on $``$. In this case, we say that $`v`$ lies above the prime number $`p`$.
### 1.2. Local fields
For each $`vM_K`$, we can form the completion $`K_v`$ (often called the local field at $`v`$) of $`K`$ with respect to $`||_v`$. We write $`_v`$ for the completion of an algebraic closure $`\overline{K}_v`$ of $`K_v`$. (The absolute value $`v`$ extends in a unique way to $`\overline{K}_v`$ and hence to $`_v`$.) The field $`_v`$ is then a complete and algebraically closed field. If $`v`$ is archimedean, then $`K_v`$ is isomorphic either to $``$ (in which case we call $`v`$ a real prime) or to $``$ (in which case we call $`v`$ a complex prime), and $`_v`$. We will henceforth avoid the notation $`K_v`$, as we will soon introduce the notation $`𝒦_v`$ to denote a completely different object in Section 2.
If $`v`$ is non-archimedean, then $`_v`$ is not locally compact, but it has other convenient properties not shared by $``$. In particular, the disk $`𝒪_v=\{c_v:|c|_v1\}`$ forms a ring, called the ring of integers, which has a unique maximal ideal $`_v=\{c_v:|c|_v<1\}`$. The quotient $`k_v=𝒪_v/_v`$ is called the residue field of $`_v`$. The natural reduction map from $`𝒪_v`$ to $`k_v`$, sending $`a𝒪`$ to $`\overline{a}=a+_vk_v`$, will be used to define good and bad reduction of a polynomial in Definition 2.1 below; but after proving a few simple Lemmas about good and bad reduction, we will not need to refer to $`𝒪_v`$, $`_v`$, or $`k_v`$ again.
### 1.3. Disks
Let $`_v`$ be a complete and algebraically closed field with absolute value $`v`$. Given $`a_v`$ and $`r>0`$, we write
$$\overline{D}(a,r)=\{x_v:|xa|_vr\}\text{and}D(a,r)=\{x_v:|xa|_v<r\}$$
for the closed and open disks, respectively, of radius $`r`$ centered at $`a`$. Note our convention that all disks have positive radius.
If $`v`$ is non-archimedean and $`U_v`$ is a disk, then the radius of $`U`$ is unique; it is the same as the diameter of the set $`U`$ viewed as a metric space. However, any point $`bU`$ is a center. That is, if $`|ba|_vr`$, then $`\overline{D}(a,r)=\overline{D}(b,r)`$, and similarly for open disks. It follows that two disks intersect if and only if one contains the other. In addition, all disks are both open and closed as topological sets; however, open disks and closed disks can still behave differently in other ways.
Still assuming that $`v`$ is non-archimedean, the set $`|_v^\times |_v`$ of absolute values actually attained by elements of $`_v^\times `$ is usually not all of $`(0,\mathrm{})`$. As a result, if $`r(0,\mathrm{})|_v^\times |_v`$, then $`D(a,r)=\overline{D}(a,r)`$ for any $`a_v`$. However, if $`r|_v^\times |_v`$, then $`D(a,r)\overline{D}(a,r)`$.
## 2. Bad Reduction and Filled Julia Sets
The following definition originally appeared in . We have modified it slightly so that “bad reduction” now means not potentially good, as opposed to not good.
###### Definition 2.1.
Let $`_v`$ be a complete, algebraically closed non-archimedean field with absolute value $`||_v`$, ring of integers $`𝒪_v=\{c_v:|c|_v1\}`$, and residue field $`k_v`$. Let $`\varphi (z)_v(z)`$ be a rational function with homogenous presentation
$$\varphi \left([x,y]\right)=[f(x,y),g(x,y)],$$
where $`f,g𝒪_v[x,y]`$ are relatively prime homogeneous polynomials of degree $`d=\mathrm{deg}\varphi `$, and at least one coefficient of $`f`$ or $`g`$ has absolute value $`1`$. We say that $`\varphi `$ has good reduction at $`v`$ if $`\overline{f}`$ and $`\overline{g}`$ have no common zeros in $`k_v\times k_v`$ besides $`(x,y)=(0,0)`$. We say that $`\varphi `$ has potentially good reduction at $`v`$ if there is some linear fractional transformation $`h\text{PGL}(2,_v)`$ such that $`h^1\varphi h`$ has good reduction. If $`\varphi `$ does not have potentially good reduction, we say it has bad reduction at $`v`$.
Naturally, for $`f[x,y]=_{i=0}^da_ix^iy^{di}`$, the reduction $`\overline{f}[x,y]`$ in Definition 2.1 means $`_{i=0}^d\overline{a}_ix^iy^{di}`$. By convention, if $`_v`$ is archimedean, we declare all rational functions in $`_v(z)`$ to have bad reduction.
In this paper, we will consider only polynomial functions $`\varphi `$ of degree at least $`2`$; that is, $`\varphi (z)=a_dz^d+\mathrm{}+a_0`$, where $`d2`$, $`a_i_v`$, and $`a_d0`$. If $`_v`$ is non-archimedean, then, it is easy to check that $`\varphi `$ has good reduction if and only if $`|a_i|_v1`$ for all $`i`$ and $`|a_d|_v=1`$. In particular, by the product formula, if $`\varphi K[z]`$ for a global field $`K`$, then there can be only finitely many primes $`vM_K`$ at which $`\varphi `$ has bad reduction.
The main focus of our investigation will be filled Julia sets, which are standard objects of study in complex dynamics. The motivating idea is that for a polynomial $`\varphi `$, any point $`x`$ of large enough absolute value will be sucked out to the attracting fixed point at $`\mathrm{}`$ under iteration; thus, all of the interesting dynamics involves points that do not escape to $`\mathrm{}`$ under iteration. Since we will be interested in both archimedean and non-archimedean fields, we state the definition here more generally.
###### Definition 2.2.
Let $`_v`$ be a complete, algebraically closed field with absolute value $`||_v`$, and let $`\varphi (z)_v[z]`$ be a polynomial of degree $`d2`$. The filled Julia set of $`\varphi `$ at $`v`$ is
$$𝒦_v=\{x_v:\{|\varphi ^n(x)|_v\}_{n1}\text{ is bounded}\}.$$
We use the notation $`𝒦_v`$ rather than $`𝒦_\varphi `$ because in this paper, we will consider the polynomial $`\varphi K[z]`$ to be fixed, and we will study its filled Julia sets at various different primes $`v`$ of $`K`$.
Non-archimedean filled Julia sets have been studied in Section 5 of , for example. It is important to note that while complex filled Julia sets are always compact, their non-archimedean counterparts are not usually compact. Fortunately, this technicality will not be an obstacle for our investigations.
We note four fundamental properties of filled Julia sets. First, $`𝒦_v`$ is invariant under $`\varphi `$; that is, $`\varphi ^1(𝒦_v)=\varphi (𝒦_v)=𝒦_v`$. Second, all the finite preperiodic points of $`\varphi `$ (that is, all the preperiodic points in $`^1(_v)`$ other than the fixed point at $`\mathrm{}`$) are contained in $`𝒦_v`$. Third, if $`U_0`$ is a disk containing $`𝒦_v`$, then
$$𝒦_v=\underset{n0}{}\varphi ^n(U_0).$$
Finally, if the polynomial $`\varphi _v[z]`$ has good reduction, then $`𝒦_v=\overline{D}(0,1)`$.
Filled Julia sets have been studied extensively in the archimedean case $`_v=`$. If $`\varphi _d(z)=z^d`$, then the (complex) filled Julia $`𝒦`$ of $`\varphi _d`$ is simply the closed unit disk $`\overline{D}(0,1)`$. Meanwhile, since the degree $`d`$ Chebyshev polynomial $`\psi _d`$ satisfies $`\psi h=h\varphi `$, where $`h(z)=z+1/z`$, it follows that the complex filled Julia set of $`\psi _d`$ is the interval $`[2,2]`$ in the real line. These two examples are misleadingly simple, however; most filled Julia sets are complicated fractal sets. For example, it is well known that for $`|c|>2`$, the filled Julia set of $`\varphi (z)=z^2+c`$ is homeomorphic to the Cantor set. For many more complex examples (often in the form of the Julia set, which is the boundary of the filled Julia set), see .
There are not as many examples of non-archimedean filled Julia sets in the literature. For the convenience of the reader, we present a few here. More examples may be found in .
###### Example 2.3.
Given $`_v`$ non-archimedean and $`d2`$, fix $`c_v`$, and consider $`\varphi (z)=z^dc^{d1}z`$. Assume for convenience that $`|d1|_v=1`$. If $`|c|_v1`$, then $`\varphi `$ has good reduction, and hence $`𝒦_v=\overline{D}(0,1)`$. Thus, we consider $`|c|_v>1`$; let $`r=|c|_v`$ and $`U=\overline{D}(0,r)`$. Note that for $`|x|_v>r`$, we have $`|\varphi (x)|_v=|x|_v^d`$, so that $`\varphi ^n(x)\mathrm{}`$. That is, $`𝒦_vU_0`$.
The set $`\varphi ^1(0)`$ consists of $`0`$ and $`d1`$ other points, all distance $`r`$ from one another. Using standard mapping properties of non-archimedean polynomials (see, for example, Section 2 of ), it is not hard to show that $`\varphi ^1(U_0)`$ consists of $`d`$ disks of radius $`r^{2d}`$, each centered at one of the points of $`\varphi ^1(0)`$. Moreover, each of these smaller disks maps one-to-one and onto $`U_0`$, and in fact $`\varphi `$ multiplies distances by a factor of $`r^{d1}=|\varphi ^{}(0)|_v`$ on each smaller disk. It follows that $`U_n=\varphi ^n(U_0)`$ is a union of $`d^n`$ disks, each of radius $`r^{1(d1)n}`$. (The sets $`U_n`$ are nested so that each disk of $`U_n`$ contains exactly $`d`$ disks of $`U_{n+1}`$, arranged so that any two are the maximal distance $`r^{1(d1)n}`$ apart.) Since the radii of the disks approach zero, it is easy to verify that $`𝒦_v=U_n`$ is homeomorphic to a Cantor set on $`d`$ intervals.
###### Example 2.4.
Given $`_v`$ non-archimedean and $`d3`$, fix $`a_v`$, and consider $`\varphi (z)=z^daz^{d1}`$. Again, assume for convenience that $`|d1|_v=1`$. If $`|a|_v1`$, then $`\varphi `$ has good reduction, and hence $`𝒦_v=\overline{D}(0,1)`$. As in the previous example, then, we consider $`|a|_v>1`$. Let $`r=|a|_v`$; once again, we have $`𝒦_vU_0=\overline{D}(0,r)`$.
This time, however, $`\varphi ^1(U_0)`$ consists of only two disks. One, $`W_1=\overline{D}(a,r^{(d1)})`$, is small and maps one-to-one onto $`U_0`$; but the other, $`W_2=\overline{D}(0,1)`$, is comparatively large, and it maps $`(d1)`$-to-$`1`$ onto $`U_0`$. Because the disk $`V=\overline{D}(0,r^{1/(d2)})`$ maps $`(d1)`$-to-1 onto itself, we can understand the shape of $`𝒦_v`$ reasonably well. In particular, $`\varphi ^2(U_0)`$ consists of $`d+2`$ disks: two inside $`W_1`$ (one is a preimage of $`W_1`$, and the other is a preimage of $`W_2`$), and $`d`$ inside $`W_2`$ ($`d1`$ are preimages of $`W_1`$, and the last is a $`(d1)`$-to-$`1`$ preimage of $`W_2`$). If we continue to take preimages, then each $`U_n=\varphi ^n(U_0)`$ will be a union of disks. To compute $`U_{n+1}`$ from $`U_n`$, observe that each disk of $`U_n`$ will will have one preimage inside $`W_1`$ and (with one exception) $`d1`$ preimages inside $`W_2`$. The one exception is the unique disk of $`U_n`$ containing $`0`$; it has only one preimage inside $`W_2`$, and that preimage maps $`(d1)`$-to-one onto it. Ultimately $`𝒦_v=U_n`$ will consist of the disk $`V`$ and all of its (infinitely many) preimages together with a vaguely Cantor-like set at which the preimages of $`V`$ accumulate.
Thus, in contrast with Example 2.3, $`𝒦_v`$ is neither a disk nor compact. In general, the filled Julia set of a polynomial of bad reduction over $`_v`$ will look something like this one. However, the dynamics can be even more complicated when there are regions on which $`\varphi `$ maps $`n`$-to-$`1`$ for some integer $`n`$ divisible by $`p`$, the characteristic of the residue field $`k_v`$.
The preceding comments and examples made frequent reference to disks $`U_0`$ containing $`𝒦_v`$. The smallest such disk will be of particular importance to us. The following Lemma shows the existence of the smallest disk and gives a partial characterization of it.
###### Lemma 2.5.
Let $`_v`$ be a complete, algebraically closed field with absolute value $`||_v`$. Let $`\varphi _v[z]`$ be a polynomial of degree $`d2`$ with lead coefficient $`a_d_v`$. Denote by $`𝒦_v`$ the filled Julia set of $`\varphi `$ in $`_v`$. Then:
* There is a unique smallest disk $`U_0_v`$ which contains $`𝒦_v`$.
* $`U_0`$ is a closed disk of some radius $`r_v^{}|_v^\times |_v`$, with $`r_v^{}|a_d|_v^{1/(d1)}`$.
* If $`||_v`$ is non-archimedean, then $`\varphi `$ has potentially good reduction if and only if $`r_v^{}=|a_d|_v^{1/(d1)}`$.
###### Proof.
Choose $`\alpha _v`$ such that $`\alpha ^{d1}=a_d`$, and let $`\psi (z)=\alpha \varphi (\alpha ^1z)`$, which is a monic polynomial with filled Julia set $`\alpha 𝒦_v`$. Given the scaling factors of $`|a_d|_v^{1/(d1)}=|\alpha |_v^1`$ in parts (b) and (c), we may assume without loss that $`a_d=1`$.
If $`_v`$ is archimedean, then $`_v`$. In that case, it is well known that $`𝒦_v`$ is a compact set in the plane. (See, for example, Lemma 9.4 of .) Since $`𝒦_v`$ is bounded, there is a unique smallest disk containing $`𝒦_v`$. (See, for example, Exercise 3 in Appendix I of .) Moreover, because $`𝒦_v`$ is compact, this disk must be closed, and so we denote it $`\overline{D}(b,r_v)`$.
It is well known that the filled Julia set of a monic polynomial over $``$ has capacity $`1`$; see, for example, Theorem 4.1 of . If $`r_v<1`$, then $`𝒦_v`$ would fit inside a disk of radius strictly smaller than $`1`$. However, the capacity of a disk in the plane is exactly its radius; as a result, the capacity of $`𝒦_v`$ would be strictly smaller than $`1`$, which is a contradiction. Therefore, $`r_v1`$, proving the Lemma in the archimedean case. (See Remark 2.6 below for an alternate proof not using capacity theory.)
If $`_v`$ is non-archimedean, then let $`b_v`$ be a fixed point of $`\varphi `$. (Such $`b`$ exists because $`\varphi (z)z`$ is a polynomial of degree $`d2`$.) Clearly $`b𝒦_v`$. By the coordinate change $`zz+b`$, we may assume that $`b=0`$. Write $`\varphi (z)=z^d+a_{d1}z^{d1}+\mathrm{}+a_1z`$. Let $`\stackrel{~}{r}_v=\mathrm{max}\{|a_i|_v^{1/(di)}:i=1,\mathrm{},d1\}`$ and $`r_v=\mathrm{max}\{\stackrel{~}{r}_v,1\}`$; note that $`r_v|_v^\times |_v`$.
If $`r_v=1`$, then $`\varphi `$ is a monic polynomial with coefficients in $`𝒪_v`$. Hence, $`\varphi `$ has good reduction; with $`U_0=\overline{D}(0,1)`$, the Lemma follows.
If $`r_v>1`$, then the Newton polygon (see Section 6.5 of or Section IV.3 of ) for the equation $`\varphi (z)=0`$ shows that there is some $`c_v`$ with $`|c|_v=r_v`$ and $`\varphi (c)=0`$. In particular, any disk containing $`𝒦_v`$ must contain $`\overline{D}(0,r_v)`$.
Moreover, if $`|z|_v>r_v`$, then the $`z^d`$ term has larger absolute value than any other term of $`\varphi (z)`$, so that $`|\varphi (z)|_v=|z|_v^d`$. By induction, $`|\varphi ^n(z)|_v=|z|_v^{d^n}`$ for all $`n1`$. It follows that $`𝒦_v\overline{D}(0,r_v)`$. By the previous paragraph, $`\overline{D}(0,r_v)`$ is the smallest disk containing $`𝒦_v`$.
It only remains to show that if $`r_v>1`$, then $`\varphi `$ cannot have potentially good reduction. However, after our coordinate changes, $`\varphi `$ is a monic polynomial without constant term. The assumption that $`r_v>1`$ means that $`|a_i|_v>1`$ for some coefficient $`a_i`$ of $`\varphi `$. Thus, by Corollary 4.6 of , $`\varphi `$ cannot have good reduction even after a change of coordinates. ∎
###### Remark 2.6.
The fact that $`r_v1`$ in the archimedean case can also be proven directly, without reference to the power of capacity theory. The following alternate argument was suggested to the author by Laura DeMarco.
Suppose that $`𝒦_v\overline{D}(a,r_v)`$ for some $`r_v<1`$; let $`s=(1+r_v)/2`$. Since $`\varphi `$ is a (monic) polynomial of degree at least $`2`$, there is some radius $`R>1`$ such that for all $`z`$ with $`|za|>R`$, we have $`|\varphi (z)a|>R`$.
Let $`A`$ denote the annulus $`\{z:s|za|R\}`$. Every point of $`A`$ is attracted to $`\mathrm{}`$ under iteration of $`\varphi `$. Since $`A`$ is compact, there is some $`n1`$ such that $`f(z)=\varphi ^n(z)a`$ has $`|f(z)|>1`$ for all $`zA`$. Note that all $`d^n`$ zeros of $`f`$ lie in $`D(a,s)`$.
Let $`g(z)=(za)^{d^n}f(z)`$, which is a polynomial of degree strictly less than $`d^n`$. However, for all $`z`$ with $`|za|=s`$, we have
$$|f(z)+g(z)|=|za|^{d^n}=s^{d^n}<1<|f(z)|.$$
By Rouché’s Theorem (noting that $`g(z)0`$ for $`|za|=s`$), $`f`$ and $`g`$ have the same number of zeros in $`D(a,s)`$, counting multiplicity. That is a contradiction; thus, $`r_v1`$.
The next two Lemmas give slightly more detailed information about the filled Julia set for a polynomial of bad reduction over a non-archimedean field.
###### Lemma 2.7.
With notation as in Lemma 2.5, suppose that $`_v`$ is non-archimedean and $`r_v^{}>|a_d|_v^{1/(d1)}`$. Then $`\varphi ^1(U_0)`$ is a disjoint union of closed disks $`D_1,\mathrm{},D_{\mathrm{}}U_0`$, where $`2\mathrm{}d`$. Moreover, there are positive integers $`d_1,\mathrm{},d_{\mathrm{}}`$ with $`d_1+\mathrm{}+d_{\mathrm{}}=d`$ such that for each $`i=1,\mathrm{},\mathrm{}`$, $`\varphi `$ maps $`D_i`$ $`d_i`$-to-$`1`$ onto $`U_0`$. That is, $`\varphi (D_i)=U_0`$, and every point $`U_0`$ has exactly $`d_i`$ pre-images in $`D_i`$, counting multiplicity.
###### Proof.
As in the previous proof, we may assume that $`\varphi (z)=z^d+a_{d1}z^{d1}+\mathrm{}+a_1z`$, that $`r_v^{}=\mathrm{max}\{|a_i|_v^{1/(di)}:1id1\}>1`$, and that $`U_0=\overline{D}(0,r_v^{})`$. As observed in that proof, $`|\varphi (z)|_v>|z|_v`$ for $`|z|_v>r_v^{}`$. Thus, $`\varphi ^1(U_0)U_0`$. Moreover, $`\varphi (U_0)U_0`$.
We now construct the disks $`D_1,\mathrm{},D_{\mathrm{}}`$ inductively. For each $`i=1,2,\mathrm{}`$, suppose we already have $`D_1,\mathrm{},D_{i1}`$, and choose $`b_i\varphi ^1(U_0)(D_1\mathrm{}D_{i1})`$. (If this is not possible, then skip to the next paragraph.) By Lemmas 2.3 and 2.6 of , there is a unique disk $`D_i`$ containing $`b_i`$ which maps onto $`U_0`$, and this disk must be closed. Since $`D_j`$ was also unique for each $`j<i`$, the new disk $`D_i`$ must be disjoint from $`D_j`$. In addition, by Lemma 2.2 of , $`\varphi `$ maps $`D_i`$ $`d_i`$-to-$`1`$ onto $`U_0`$, for some integer $`d_i1`$.
This process must stop with $`\mathrm{}d`$, because $`\varphi ^1(0)`$ consists of exactly $`d`$ points, counting multiplicity, and since each $`d_i1`$, at least one must be contained in each $`D_i`$. In fact, counting pre-images of $`0`$ also shows that $`d_1+\mathrm{}+d_{\mathrm{}}=d`$.
Finally, suppose that $`\mathrm{}=1`$; that is, $`\varphi ^1(U_0)`$ is a single disk $`D_1U_0`$. However, $`𝒦_v=\varphi ^1(𝒦_v)D_1`$, contradicting the assumption that $`U_0`$ is the smallest disk containing $`𝒦_v`$. Thus, we must have $`\mathrm{}2`$. ∎
###### Lemma 2.8.
Let $`K`$ be a field with a discrete valuation $`v`$, and let $`\pi _vK`$ be a uniformizer at $`v`$. Let $`_v`$ be the completion of an algebraic closure of $`K`$.
Let $`\varphi (z)=a_dz^d+\mathrm{}+a_0K[z]`$ be a polynomial of degree $`d2`$. Denote by $`𝒦_v`$ the filled Julia set of $`\varphi `$ in $`_v`$, and let $`r_v^{}>0`$ be the radius of the smallest disk in $`_v`$ containing $`𝒦_v`$. Suppose that $`r_v^{}>|a_d|_v^{1/(d1)}`$.
If $`𝒦_vK\mathrm{}`$, then
$$|a_d|_v^{1/(d1)}r_v^{}\{\begin{array}{cc}|\pi _v|_v^1\hfill & \text{ if }d=2,\hfill \\ |\pi _v|_v^{1/[(d1)(d2)]}\hfill & \text{ if }d3.\hfill \end{array}$$
###### Proof.
Given $`b𝒦_vK`$, we may replace $`\varphi `$ by $`\varphi (z+b)bK[z]`$, which is a polynomial with the same degree and lead coefficient as $`\varphi `$, but with filled Julia set translated by $`b`$. In particular, the radius $`r_v^{}`$ is preserved; so we may assume without loss that $`0𝒦_v`$.
As in the proof of Lemma 2.5, choose $`\alpha _v`$ such that $`\alpha ^{d1}=a_d`$, and let
$$\psi (z)=\alpha \varphi (\alpha ^1z)=z^d+\underset{i=0}{\overset{d1}{}}\alpha ^{1i}a_iz^i.$$
Then $`\psi `$ is a monic polynomial with filled Julia set $`𝒦_v^{}=\alpha 𝒦_v`$; however, $`\psi `$ may not be defined over $`K`$. Still, the radius $`r_v`$ of the smallest disk containing $`𝒦_v^{}`$ satisfies $`r_v>1`$, by hypothesis.
Let $`j`$ be the largest index between $`0`$ and $`d1`$ that maximizes $`\lambda _j=|\alpha ^{1j}a_j|_v^{1/(dj)}`$. Note that $`\lambda _j>1`$; for if not, then $`\psi `$ has good reduction, contradicting Lemma 2.5.c. The Newton polygon for the equation $`\psi (z)=0`$ shows that there is some $`\beta _v`$ with $`\psi (\beta )=0`$ and $`|\beta |_v=\lambda _j`$. We have $`0,\beta 𝒦_v^{}`$; hence, $`r_v\lambda _j`$.
If $`j=0`$, then a simple induction shows that $`|\psi ^n(0)|_v=|\alpha a_0|_v^{d^{n1}}`$ for $`n1`$. Since $`|\alpha a_0|_v>1`$, this contradicts the hypothesis that $`0𝒦_v^{}`$.
Thus, $`1jd1`$, and we write $`|a_d|_v=|\pi _v|_v^{e_1}`$ and $`|a_j|_v=|\pi _v|_v^{e_2}`$; note that $`e_1,e_2`$. Our assumptions say that
$$r_v\lambda _j=|\alpha ^{1j}a_j|_v^{1/(dj)}=|\pi _v|_v^f>1,\text{where}f=\frac{1}{dj}\left(\frac{(1j)e_1}{d1}+e_2\right)<0.$$
If $`j=1`$, then $`f=e_2/(d1)1/(d1)`$, which proves the Lemma for $`d=2`$. If $`2jd1`$, then $`f1/[(d1)(dj)]1/[(d1)(d2)]`$, and we are done. ∎
###### Remark 2.9.
The bounds of Lemma 2.8 are sharp. Indeed, if $`d=2`$, then the polynomial $`\varphi (z)=z^2\pi _v^1z`$ has $`0,\pi _v^1𝒦_v`$. Because $`\varphi ^n(z)\mathrm{}`$ for $`|z|_v>|\pi _v|_v^1`$, the smallest disk containing $`𝒦_v`$ is $`\overline{D}(0,|\pi _v|_v^1)`$, so that $`r_v=|\pi _v|_v^1`$.
Similarly, if $`d3`$, then the polynomial $`\varphi (z)=\pi _v^dz^d\pi _vz^2`$ has $`0𝒦_v`$. Choosing a $`(d1)`$-st root of $`\pi _v`$, we get an associated monic conjugate $`\psi (z)=z^d\pi _v^{1/(d1)}z^2`$, from which it is easy to compute $`|\pi _v^d|_v^{1/(d1)}r_v^{}=r_v=|\pi _v|_v^{1/[(d1)(d2)]}`$.
## 3. Elementary Computations
We will now define and bound certain integer quantities that will appear as exponents in the rest of the paper. The reader is encouraged to read the statements of Definition 3.1, Lemma 3.4, and Lemma 3.5 but to skip the proofs, which are tedious but completely elementary, until after seeing their use in Theorem 7.1.
We will write $`\mathrm{log}_dx`$ to denote the logarithm of $`x`$ to base $`d`$.
###### Definition 3.1.
Let $`N0`$ and $`d2`$ be integers. We define $`E(N,d)`$ to be twice the sum of all base-$`d`$ coefficients of all integers from $`0`$ to $`N1`$. That is,
$$E(N,d)=2\underset{j=0}{\overset{N1}{}}e(j,d),\text{where}e(\underset{i=0}{\overset{M}{}}c_id^i,d)=\underset{i=0}{\overset{M}{}}c_i,$$
for $`c_i\{0,1,\mathrm{},d1\}`$.
Moreover, if $`m`$ is an integer satisfying $`1md`$, we may write $`N=c_0+mk`$ for unique integers $`c_0\{0,1,\mathrm{},m1\}`$ and $`k0`$. We then define
$$e(N,m,d)=c_0+e(k,d)(dm)k\text{and}f(N,m,d)=c_0+e(k,d),$$
and
$$E(N,m,d)=2\underset{j=0}{\overset{N1}{}}e(j,m,d)\text{and}F(N,m,d)=2\underset{j=0}{\overset{N1}{}}f(j,m,d).$$
We declare $`E(N,d)=E(N,m,d)=F(N,m,d)=0`$ for $`N1`$. Clearly, $`E(N,d)`$ and $`F(N,m,d)`$ are always positive for $`N1`$; but for $`N`$ large and $`m<d`$, $`E(N,m,d)`$ is negative. Note that $`e(N,d,d)=f(N,d,d)=f(N,1,d)=e(N,d)`$, and therefore
(3)
$$E(N,d,d)=F(N,d,d)=F(N,1,d)=E(N,d).$$
We will need the following two auxiliary Lemmas.
###### Lemma 3.2.
Let $`N,m,d`$ be integers satisfying $`N1`$, $`d2`$, and $`1md`$. Write $`N=c+mk`$ with $`0cm1`$ and $`k0`$. Then:
* $`F(N,m,d)=(mc)E(k,d)+cE(k+1,d)+(m1)Nc(mc)`$.
* $`E(N,m,d)=F(N,m,d){\displaystyle \frac{(dm)}{m}}[N^2mN+c(mc)]`$.
* If $`Nm`$, then $`E(N,m,d)=F(N,m,d)=N(N1)`$.
###### Proof.
Writing an arbitrary integer $`j0`$ as $`j=i+m\mathrm{}`$ for $`0im1`$, we compute
$`F(N,m,d)`$ $`=2{\displaystyle \underset{j=0}{\overset{N1}{}}}f(j,m,d)=2{\displaystyle \underset{i=0}{\overset{c1}{}}}{\displaystyle \underset{\mathrm{}=0}{\overset{k}{}}}f(i+m\mathrm{},m,d)+2{\displaystyle \underset{i=c}{\overset{m1}{}}}{\displaystyle \underset{\mathrm{}=0}{\overset{k1}{}}}f(i+m\mathrm{},m,d)`$
$`=2{\displaystyle \underset{i=0}{\overset{c1}{}}}{\displaystyle \underset{\mathrm{}=0}{\overset{k}{}}}(i+e(\mathrm{},d))+2{\displaystyle \underset{i=c}{\overset{m1}{}}}{\displaystyle \underset{\mathrm{}=0}{\overset{k1}{}}}(i+e(\mathrm{},d))`$
$`={\displaystyle \underset{i=0}{\overset{c1}{}}}\left[2(k+1)i+E(k+1,d)\right]+{\displaystyle \underset{i=c}{\overset{m1}{}}}\left[2ki+E(k,d)\right]`$
$`=cE(k+1,d)+(mc)E(k,d)+(k+1)c(c1)+km(m1)kc(c1).`$
Part (a) of the Lemma now follows by rewriting the last three terms as
$$c(c1)+mk(m1)=c(cm)+(c+mk)(m1)=(m1)Nc(mc).$$
Next, we compute
$`E(N,m,d)`$ $`=2{\displaystyle \underset{j=0}{\overset{N1}{}}}e(j,m,d)=2{\displaystyle \underset{j=0}{\overset{N1}{}}}f(j,m,d)2(dm)\left[m{\displaystyle \underset{\mathrm{}=0}{\overset{k1}{}}}\mathrm{}+{\displaystyle \underset{j=0}{\overset{c1}{}}}k\right]`$
$`=F(N,m,d)k(dm)[m(k1)+2c]`$
Writing $`k=(Nc)/m`$, the last term becomes
$$\frac{(Nc)}{m}(dm)(N+cm)=\frac{(dm)}{m}[N^2mN+c(mc)],$$
proving part (b). Finally, part (c) is immediate from the fact that $`e(j,m,d)=f(j,m,d)=j`$ for $`0jm1`$. ∎
###### Lemma 3.3.
Let $`N,m,d`$ be integers satisfying $`N1`$, $`d2`$, and $`1md`$. Write $`N=c+mk`$ with $`0cm1`$ and $`k0`$. Then:
* $`(mc)\mathrm{log}_d\left({\displaystyle \frac{mk}{N}}\right)+c\mathrm{log}_d\left({\displaystyle \frac{mk+m}{N}}\right)0`$.
* If $`Nd`$, then $`(d1)\mathrm{log}_d\left({\displaystyle \frac{mk+m}{N}}\right)(mc)0`$.
###### Proof.
The function $`\mathrm{log}_d(x)`$ is of course concave down. Letting $`x_1=mk/N`$ and $`x_2=(mk+m)/N`$, then, we have $`x_11<x_2`$, and therefore $`\mathrm{log}_d(1)L(1)`$, where
$$L(x)=\frac{1}{x_2x_1}\left[(x_2x)\mathrm{log}_d(x_1)+(xx_1)\mathrm{log}_d(x_2)\right]$$
is the line through $`(x_1,\mathrm{log}_d(x_1))`$ and $`(x_2,\mathrm{log}_d(x_2))`$. That is,
$$0\frac{1}{m}\left[(mc)\mathrm{log}_d\left(\frac{mk}{N}\right)+c\mathrm{log}_d\left(\frac{mk+m}{N}\right)\right],$$
proving part (a). For part (b), we have
$$(d1)\mathrm{log}_d\left(\frac{mk+m}{N}\right)=\frac{(d1)}{\mathrm{log}d}\mathrm{log}\left(1+\frac{mc}{N}\right)\frac{(d1)}{\mathrm{log}d}\frac{(mc)}{N}.$$
However, $`\mathrm{log}d=\mathrm{log}[1(d1)/d](d1)/d`$, and since $`Nd`$,
$$(d1)\mathrm{log}_d\left(\frac{mk+m}{N}\right)(d1)\frac{d}{d1}\frac{mc}{N}=\frac{d}{N}(mc)(mc).\mathit{}$$
###### Lemma 3.4.
Let $`N,m,d`$ be integers satisfying $`N1`$, $`d2`$, and $`1md1`$. Then:
* $`E(N,d)(d1)N\mathrm{log}_dN`$, with equality if $`N`$ is a power of $`d`$.
* $`E(N,m,d)(d1)N\left[\mathrm{log}_dN+1\mathrm{log}_dm{\displaystyle \frac{(dm)}{m(d1)}}N\right]`$, with equality if $`N/m`$ is a power of $`d`$.
* $`F(N,m,d)(d1)N\mathrm{log}_dN`$.
* For $`Nm`$, $`F(N,m,d)(d1)N\left[\mathrm{log}_dN\mathrm{log}_dm+{\displaystyle \frac{m1}{d1}}\right]`$, with equality if $`N/m`$ is a power of $`d`$.
###### Proof.
Fix $`d2`$. If $`N=1`$, then both sides of part (a) are clearly zero. If $`2Nd`$, then $`E(N,d)=N(N1)`$ by Lemma 3.2.c (with $`m=d`$) and equation (3). Because $`(\mathrm{log}x)/(x1)`$ is a decreasing function of the real variable $`x>1`$, we have $`(\mathrm{log}d)/(d1)(\mathrm{log}N)/(N1)`$, with equality for $`N=d`$. Part (a) then follows for $`1Nd`$.
For $`Nd+1`$, we proceed by induction on $`N`$, assuming part (a) holds for all positive integers up to $`N1`$. Write $`N=c+dk`$, where $`0cd1`$, so that $`1kN2`$. By Lemma 3.2.a (with $`m=d`$) and equation (3), we have
$`E(N,d)`$ $`=(dc)E(k,d)+cE(k+1,d)+(d1)Nc(dc)`$
$`(dc)(d1)k\mathrm{log}_dk+c(d1)(k+1)\mathrm{log}_d(k+1)+(d1)Nc(dc)`$
$`=(dc)(d1)k\mathrm{log}_d(dk)+c(d1)(k+1)\mathrm{log}_d(dk+d)c(dc)`$
where the final equality is because $`N=(dc)k+c(k+1)`$, and the inequality (which is equality if $`N`$ is a power of $`d`$) is by the inductive hypothesis, since $`k,k+1N1`$. More generally, adding and subtracting $`(d1)N\mathrm{log}_dN`$, we have
$`E(N,d)`$ $`(d1)N\mathrm{log}_dN+(dc)(d1)k\mathrm{log}_d\left({\displaystyle \frac{dk}{N}}\right)`$
$`+c(d1)(k+1)\mathrm{log}_d\left({\displaystyle \frac{dk+d}{N}}\right)c(dc)`$
$`=(d1)N\mathrm{log}_dN+c\left[(d1)\mathrm{log}_d\left({\displaystyle \frac{dk+d}{N}}\right)(dc)\right].`$
$`+(d1)k\left[(dc)\mathrm{log}_d\left({\displaystyle \frac{dk}{N}}\right)+c\mathrm{log}_d\left({\displaystyle \frac{dk+d}{N}}\right)\right]`$
By Lemma 3.3 with $`m=d`$, the quantities in brackets are nonpositive, and part (a) follows.
If $`m=1`$, then parts (c–d) are immediate from part (a) and equation (3). Moreover, by Lemma 3.2.a–b (with $`m=1`$) and part (a),
$$E(N,1,d)=E(N,d)(d1)N(N1)(d1)N[\mathrm{log}_dN+1N],$$
with equality if $`N`$ is a power of $`d`$. This is exactly part (b) for $`m=1`$. Thus, we may assume for the remainder of the proof that $`2md`$.
We now turn to part (d). If $`N=m`$, then by Lemma 3.2.c, we have $`F(m,m,d)m(m1)`$, which exactly equals the desired right hand side. For $`Nm+1`$, write $`N=c+mk`$, where $`k1`$ and $`0cm1`$. By Lemma 3.2.a,
(4)
$$F(N,m,d)=(mc)E(k,d)+cE(k+1,d)+(m1)Nc(mc).$$
If $`m+1Nd1`$, then $`kd1`$, so that by Lemma 3.2.c, equation (4) becomes
$`F(N,m,d)`$ $`=(mc)k(k1)+ck(k+1)+(m1)Nc(mc)`$
$`=mk^2mk+2ck+(m1)Nc(mc)=(m+k1)N(mc)(k+c)`$
$`=m^1\left[(N+m^2cm)N(mc)(Nc+cm)\right]`$
$`=m^1\left[(N+m^22m)Nc(mc)(m1)\right]N\left[(N/m)+m2\right],`$
where we have substituted $`k=(Nc)/m`$ along the way. Thus, we must show
$$N\left[(N/m)+m2\right]N\left[(d1)\mathrm{log}_d(N/m)+m1\right].$$
Equivalently, we must show
$$\frac{\mathrm{log}d}{d1}\frac{\mathrm{log}(N/m)}{(N/m)1},$$
which is true because $`(\mathrm{log}x)/(x1)`$ is a decreasing function of $`x>1`$, and $`1<N/m<d`$.
If $`Nd`$ in part (d), then $`k1`$. Applying part (a) to equation (4), we obtain
(5)
$$\begin{array}{c}F(N,m,d)(mc)(d1)k\mathrm{log}_dk+c(d1)(k+1)\mathrm{log}_d(k+1)\hfill \\ \hfill +(m1)Nc(mc),\end{array}$$
with equality if $`c=0`$ and $`k`$ is a power of $`d`$, whence we immediately obtain the statement of the Lemma for $`N=md^i`$. More generally, (5) becomes
$`F(N,m,d)`$ $`(d1)N\left(\mathrm{log}_d{\displaystyle \frac{N}{m}}+{\displaystyle \frac{m1}{d1}}\right)+(mc)(d1)k\mathrm{log}_d\left({\displaystyle \frac{mk}{N}}\right)`$
$`+c(d1)(k+1)\mathrm{log}_d\left({\displaystyle \frac{mk+m}{N}}\right)c(mc)`$
$`=(d1)N\left(\mathrm{log}_d{\displaystyle \frac{N}{m}}+{\displaystyle \frac{m1}{d1}}\right)+c\left[(d1)\mathrm{log}_d\left({\displaystyle \frac{mk+m}{N}}\right)(mc)\right]`$
$`+(d1)k\left[(mc)\mathrm{log}_d\left({\displaystyle \frac{mk}{N}}\right)+c\mathrm{log}_d\left({\displaystyle \frac{mk+m}{N}}\right)\right].`$
Part (d) now follows from Lemma 3.3, as before.
For part (c), if $`1Nm`$, then by part (a) and Lemma 3.2.c,
$$F(N,m,d)=N(N1)=E(N,d)(d1)N\mathrm{log}_dN,$$
as desired. The remaining case, that $`Nm`$, will follow from part (d) provided
$$m1(d1)\mathrm{log}_dm.$$
However, this is the same as showing that $`\mathrm{log}d/(d1)\mathrm{log}m/(m1)`$, which once again follows from the fact that $`(\mathrm{log}x)/(x1)`$ is decreasing for $`x>1`$.
Last, we turn to part (b). If $`1Nm`$, then $`E(N,m,d)=N(N1)`$ by Lemma 3.2.c. Thus, we wish to show that
$$N1(d1)\left(1+\mathrm{log}_d\frac{N}{m}\right)(dm)\frac{N}{m},$$
which is to say
$$\frac{dN}{m}1(d1)\mathrm{log}_d\left(\frac{dN}{m}\right),$$
with equality when $`N=m`$. Yet again, this inequality follows immediately from the facts that $`(\mathrm{log}x)/(x1)`$ is decreasing for $`x>1`$ and that $`1dN/md`$.
It only remains to consider $`Nm+1`$. Writing $`N=c+mk`$, where $`k1`$ and $`0cm1`$, and invoking Lemma 3.2.b, we have
$`E(N,m,d)`$ $`=F(N,m,d){\displaystyle \frac{(dm)}{m}}[N^2mN+c(mc)]`$
$`F(N,m,d){\displaystyle \frac{(dm)}{m}}N^2+(dm)N,`$
with equality for $`c=0`$. By part (d), we obtain
$`E(N,m,d)`$ $`(d1)N\left[\mathrm{log}_dN\mathrm{log}_dm+{\displaystyle \frac{m1}{d1}}\right]{\displaystyle \frac{(dm)}{m}}N^2+(dm)N`$
$`=(d1)N\left[\mathrm{log}_dN\mathrm{log}_dm+1{\displaystyle \frac{dm}{m(d1)}}N\right],`$
with equality if $`N`$ is of the form $`N=md^i`$. ∎
Besides the preceding integer quantities and their bounds, we will need the following bound involving a certain family of real-valued functions.
###### Lemma 3.5.
Let $`d2`$ be an integer, and let $`A,B,t`$ be positive real numbers such that
$$(d1)Ad^{B1}\text{and}t1.$$
Define $`\eta :(0,\mathrm{})`$ by
$$\eta (x)=t\mathrm{log}_dxAx+B.$$
Set the real number $`M(A,B,t)`$ to be
$$M(A,B,t)=\frac{t}{A}\left(\mathrm{log}_dt+\mathrm{log}_d(\mathrm{max}\{1,\mathrm{log}_dt\})+3\right).$$
Then $`\eta (x)<0`$ for all $`xM(A,B,t)`$.
###### Proof.
By differentiating, we see that $`\eta `$ is decreasing for $`xt/(A\mathrm{log}d)`$, and hence for $`xM(A,B,t)`$. Thus, it suffices to show that $`\eta (M(A,B,t))<0`$.
First, suppose that $`t<d`$. Then
$`\eta (M(A,B,t))`$ $`=t\mathrm{log}_dt+t\mathrm{log}_d\left[A^1(\mathrm{log}_dt+3)\right]t\mathrm{log}_dt3t+B`$
$`=t\mathrm{log}_d\left[A^1d^{B3}(\mathrm{log}_dt+3)\right]B(t1)t\mathrm{log}_d\left[(d1)(\mathrm{log}_dt+3)/d^2\right],`$
where the inequality is because $`A^1d^{B1}(d1)`$ and $`B>0`$, by hypothesis. Since $`t<d`$, the quantity inside square brackets is strictly less than $`4(d1)/d^21`$. Thus, $`\eta (M(A,B,t))<t\mathrm{log}_d(1)=0`$, and we are done.
Second, if $`td`$, then by a similar computation,
$`\eta (M(A,B,t))`$ $`=t\mathrm{log}_d\left[A^1d^{B3}u^1(u+\mathrm{log}_du+3)\right]B(t1)`$
$`<t\mathrm{log}_d\left[(d1)(u+\mathrm{log}_du+3)/(d^2u)\right]`$
where $`u=\mathrm{log}_dt`$. Writing $`H(u)=(d1)(u+\mathrm{log}_du+3)/(d^2u)`$, it suffices to show that $`H(u)1`$ for $`u1`$. Differentiating, it is easy to see that $`H`$ is decreasing for such $`u`$. Since $`H(1)=4(d2)/d^21`$, we are done. ∎
## 4. Transfinite Diameters and Bad Primes
Given a metric space $`X`$ and an integer $`N2`$, the $`N^{\text{th}}`$ diameter of $`X`$ is defined to be
$$𝐝_N(X)=\underset{x_1,\mathrm{},x_NX}{sup}\underset{ij}{}d_X(x_i,x_j)^{1/[N(N1)]},$$
which measures the maximal average distance between any two of $`N`$ points in $`X`$. (See , for example, for a computation of the $`N^{\text{th}}`$ diameter of the interval $`[0,1]`$.) This quantity is usually used to define the transfinite diameter of $`X`$,
$$𝐝(X)=\underset{N\mathrm{}}{lim}𝐝_N(X),$$
which converges because $`\{𝐝_N(X)\}_{N2}`$ is a decreasing sequence. If $`X`$ is a nice enough (e.g., compact) subset of a valued field, then the transfinite diameter coincides with the Chebyshev constant and the logarithmic capacity of $`X`$; see Section 5.4 of , or Chapters 3 and 4 of . Baker and Hsia used this equality in to compute the transfinite diameter of filled Julia sets of polynomials, even when those sets were not compact. (Their result of $`|a_d|_v^{1/(d1)}`$, where $`d`$ is the degree and $`a_d`$ the lead coefficient of the polynomial, was already well known for $`_v=`$.) See for more on transfinite diameters and capacities in $`_v`$.
However, in this paper we will be interested in the $`N^{\text{th}}`$ diameters $`𝐝_N(X)`$ themselves, rather than the transfinite diameter. In particular, the following Lemma contains our main bound for $`𝐝_N(𝒦_v)^{N(N1)}`$, where $`𝒦_v`$ is the filled Julia set of a polynomial $`\varphi _v[z]`$. The proof uses an estimate involving van der Monde determinants similar to a bound that appears in the proof of Lemme 5.4.2 in .
###### Lemma 4.1.
Let $`_v`$ be a complete, algebraically closed field with absolute value $`||_v`$. Let $`\varphi _v[z]`$ be a polynomial of degree $`d2`$ with lead coefficient $`a_d_v`$. Denote by $`𝒦_v`$ the filled Julia set of $`\varphi `$ in $`_v`$, and let $`r_v^{}`$ be the radius of the smallest disk that contains $`𝒦_v`$. Set $`r_v=|a_d|_v^{1/(d1)}r_v^{}`$.
Then for any integer $`N2`$ and any set $`\{x_1,\mathrm{},x_N\}𝒦_v`$ of $`N`$ points in $`𝒦_v`$,
$$\underset{ij}{}|x_ix_j|_v|a_d|_v^{N(N1)/(d1)}\mathrm{max}\{1,|N|_v^N\}r_v^{E(N,d)},$$
where $`E(N,d)`$ is twice the sum of all base-$`d`$ coefficients of all integers from $`0`$ to $`N1`$, as in Definition 3.1.
###### Proof.
Choose $`\alpha _v`$ such that $`\alpha ^{d1}=a_d`$, and let $`\psi (z)=\alpha \varphi (\alpha ^1z)`$. Then $`\psi `$ is a monic polynomial with filled Julia set $`𝒦_v^{}=\alpha 𝒦_v`$, and the smallest disk containing $`𝒦_v^{}`$ has radius $`r_v`$. If the Lemma holds for $`\psi `$, then for $`x_1,\mathrm{},x_N𝒦_v`$, we have $`\alpha x_i𝒦_v^{}`$, and therefore
$$\underset{ij}{}|x_ix_j|_v=|\alpha |_v^{N(N1)}\underset{ij}{}|\alpha x_i\alpha x_j|_v|\alpha |_v^{N(N1)}\mathrm{max}\{1,|N|_v^N\}r_v^{E(N,d)},$$
as desired. Thus, it suffices to prove the Lemma in the case that $`\varphi `$ is monic.
We will now construct a sequence $`\{f_j\}_{j=1}^{\mathrm{}}`$ of monic polynomials over $`_v`$ such that each $`f_j`$ has degree $`j`$ and such that $`|f_j(x)|_v`$ is not especially large for any $`x𝒦_v`$.
First, let $`\overline{D}(a,r_v)`$ be the smallest disk containing $`𝒦_v`$, where $`a_v`$ and $`r_v`$ is as in the statement of the Lemma. For any integer $`j0`$ written in base-$`d`$ notation as
$$j=c_0+c_1d+c_2d^2+\mathrm{}+c_Md^M,$$
with $`c_i\{0,1,\mathrm{},d1\}`$, define
$$f_j(z)=\underset{i=0}{\overset{M}{}}[\varphi ^i(z)a]^{c_i}.$$
Clearly, $`f_j`$ is monic of degree $`j`$. Moreover, for $`x𝒦_v`$, we have $`\varphi ^i(x)𝒦_v`$, and therefore
$$|f_j(x)|_v\underset{i=0}{\overset{M}{}}r_v^{c_i}=r_v^{e(j,d)},$$
where $`e(j,d)`$ is as in Definition 3.1.
Given $`x_1,\mathrm{},x_N𝒦_v`$, denote by $`V(x_1,\mathrm{},x_N)`$ the corresponding van der Monde matrix (i.e., the $`N\times N`$ matrix with $`(i,j)`$ entry $`x_i^{j1}`$). Recall that
$$\underset{ij}{}|x_ix_j|_v=|detV(x_1,\mathrm{},x_N)|_v^2.$$
Because $`f_{N1}`$ is monic, we may replace the last column of the matrix by a column with entry $`f_{N1}(x_i)`$ in the $`i`$th row, without changing the determinant. We may then replace the second to last column by a column with entry $`f_{N2}(x_i)`$ in the $`i`$th row, and so on. Thus, if we denote by $`A(x_1,\mathrm{},x_N)`$ the matrix with $`(i,j)`$ entry $`f_{j1}(x_i)`$, then
$$detV(x_1,\mathrm{},x_N)=detA(x_1,\mathrm{},x_N).$$
If $`_v=`$ is archimedean, then by Hadamard’s inequality applied to the columns of $`A`$,
$$|detA(x_1,\mathrm{},x_N)|^2\underset{j=0}{\overset{N1}{}}\left(|f_j(x_1)|^2+\mathrm{}+|f_j(x_N)|^2\right)\underset{j=0}{\overset{N1}{}}Nr_v^{2e(j,d)}=N^Nr_v^{E(N,d)}.$$
Similarly, if $`_v`$ is non-archimedean, then by the non-archimedean version of Hadamard’s inequality (see, for example, , Preuve du Lemme 5.3.4), we have
$$|detA(x_1,\mathrm{},x_N)|^2\underset{j=0}{\overset{N1}{}}\underset{i=1,\mathrm{},N}{\mathrm{max}}|f_j(x_i)|^2\underset{j=0}{\overset{N1}{}}r_v^{2e(j,d)}=r_v^{E(N,d)}.\mathit{}$$
###### Remark 4.2.
We can recover the Baker and Hsia bound $`𝐝(𝒦_v)|a_d|_v^{1/(d1)}`$ immediately from Lemmas 4.1 and 3.4.a. (The opposite inequality is more subtle, however.)
###### Remark 4.3.
There are many cases for which the bound of Lemma 4.1 is sharp. In particular, for non-archimedean $`v`$, degree $`d2`$ with $`|d1|_v=1`$, and $`c_v`$ with $`|c|>1`$, recall that the function $`\varphi (z)=z^dc^{d1}z`$ of Example 2.3 has $`𝒦_v`$ homeomorphic to a Cantor set on $`d`$ pieces. For arbitrary $`N2`$, one can distribute $`N`$ points in $`𝒦_v`$ in the following way. Write $`N=_{i=0}^Mc_id^i`$, and put $`c_M`$ points in each of the $`d^M`$ pieces at level $`M`$, maximally far apart in each piece; then put $`c_{M1}`$ in each of the $`d^{M1}`$ pieces at level $`M1`$, each as far as possible from the existing points; and so on. Keeping track of the radii of the disks at each level, one can show that $`_{ij}|x_ix_j|_v=r_v^{E(N,d)}`$ exactly.
In many other cases, however, the bound is not quite sharp, though it appears to be approximately the right order of magnitude. In the archimedean case, of course, the Hadamard inequality introduces some error. Still, the greater factor seems to be the choice of the monic polynomial $`f_j`$. When $`j`$ is a power of $`d`$, computations suggest that our choice of $`f_j`$ is very close to sharp, if not actually sharp. However, when $`j`$ is not a power of $`d`$, our construction of $`f_j`$ as a product of smaller factors is in general not optimal, even in the non-archimedean setting. For example, if $`\varphi (z)=z^3az^2`$ is the map of Example 2.4 (non-archimedean, with $`d=3`$, $`|a|_v>1`$, and $`|2|_v=1`$), then the function $`f_6(z)=(\varphi (z))^2`$ of the proof has $`|f_6(z)|_v`$ growing as large as $`r^2`$ on $`𝒦_v`$; but the function $`\stackrel{~}{f}_6(z)=(\varphi (z))(\varphi (z)a)`$ has $`|\stackrel{~}{f}_6(z)|_vr`$. Ultimately, while the exponent $`E(N,3)`$ of Lemma 4.1 is essentially $`2N\mathrm{log}_3N`$, the actual exponent for this $`\varphi `$ should be something more like $`(4/3)N\mathrm{log}_3N`$.
In the archimedean case, the Chebyshev polynomials $`\{\psi _j\}_{j1}`$ provide an even stronger example of this phenomenon. More precisely, if $`_v=`$ and $`\varphi (z)=\psi _2(z)=z^22`$, then $`𝒦_v`$ is simply the interval $`[2,2]`$ in the real line. For $`j1`$, the $`j^{\text{th}}`$ Chebyshev polynomial $`\psi _j`$ has $`|\psi _j|2`$ on $`𝒦_v`$, as compared with the proof’s bound of $`2^{c_0+c_1+\mathrm{}}`$ for $`|f_j|`$.
In general, however, knowing nothing about the polynomial other than its degree and the radius $`r_v`$, we cannot substantially improve on Lemma 4.1.
## 5. A Partition of the Filled Julia Set: Non-archimedean Case
The key to the Main Theorem, as described in the introduction, is to divide the filled Julia set at a particular bad prime into two smaller pieces $`X_1`$ and $`X_2`$. As a result, the product $`_{ij}|x_ix_i|_v`$, when restricted to $`\{x_i\}X_k`$ (for fixed $`k=1,2`$), will be substantially smaller than the bound of Lemma 4.1. We begin with non-archimedean primes.
###### Lemma 5.1.
Let $`_v`$ be a complete, algebraically closed field with non-archimedean absolute value $`||_v`$. Let $`\varphi _v[z]`$ be a polynomial of degree $`d2`$ with lead coefficient $`a_d_v`$. Denote by $`𝒦_v`$ the filled Julia set of $`\varphi `$ in $`_v`$, and let $`r_v^{}`$ be the radius of the smallest disk $`U_0`$ that contains $`𝒦_v`$. Set $`r_v=|a_d|_v^{1/(d1)}r_v^{}`$, and suppose that $`r_v>1`$.
Then there are disjoint sets $`X_1,X_2𝒦_v`$ and positive integers $`m_1,m_2`$ with the properties that $`X_1X_2=𝒦_v`$, that $`m_1+m_2=d`$, that for $`k=1,2`$, $`\varphi :X_k𝒦_v`$ is $`m_k`$-to-$`1`$, and that for $`k=1,2`$, for any integer $`N2`$, and for any set $`\{x_1,\mathrm{},x_N\}X_k`$ of $`N`$ points in $`X_k`$,
$$\underset{ij}{}|x_ix_j|_v|a_d|_v^{N(N1)/(d1)}r_v^{E(N,m_k,d)},$$
where $`E(N,m_k,d)`$ is as in Definition 3.1
###### Proof.
As in the proof of Lemma 4.1, we may assume that $`\varphi `$ is monic.
By Lemma 2.5, $`U_0`$ is a closed disk of radius $`r_v|_v^\times |_v`$. We may write $`U_0=\overline{D}(a,r_v)`$ for some point $`a𝒦_v`$, since $`𝒦_v`$ is nonempty, and since any point of a non-archimedean disk is a center. Pick $`b\varphi ^1(a)`$. Note that $`b𝒦_vU_0`$.
Write $`U_1=\varphi ^1(U_0)`$. By Lemma 2.7, $`U_1=D_1\mathrm{}D_{\mathrm{}}`$ for some disjoint closed disks $`\{D_i\}`$, with $`2\mathrm{}d`$. Moreover, $`\varphi :D_iU_0`$ maps $`d_i`$-to-one for some positive integers $`\{d_i\}`$ with $`d_1+\mathrm{}+d_{\mathrm{}}=d`$. Define
$$W_1=\{xU_1:|xb|_v<r_v\},\text{and}W_2=U_1W_1,$$
so that $`W_1W_2=\mathrm{}`$ and $`W_1W_2=U_1`$. If $`W_2=\mathrm{}`$, then $`𝒦_vD(b,r_v)U_0`$, contradicting the minimality of $`U_0`$. (The second inclusion is strict because $`r_v|_v^\times |_v`$.) Thus, since $`bW_1`$, both $`W_1`$ and $`W_2`$ are nonempty.
Furthermore, $`W_1`$ and $`W_2`$ are both finite unions of disks $`D_i`$ above. Hence, there are integers $`m_1,m_21`$ so that each $`W_k`$ maps $`m_k`$-to-one onto $`U_0`$, with $`m_1+m_2=d`$. Let $`X_k=W_k𝒦_v`$ for $`k=1,2`$. Since $`\varphi ^1(𝒦_v)=𝒦_v`$, $`\varphi `$ must map $`X_k`$ $`m_k`$-to-one onto $`𝒦_v`$.
For any integer $`i1`$, observe that the polynomial $`\varphi ^i(z)a`$ is monic of degree $`d^n`$. Moreover, since the equation $`\varphi ^{i1}(z)=a`$ has exactly $`d^{i1}`$ roots (counting multiplicity), all of which lie in $`U_0`$, it follows that $`\varphi ^i(z)=a`$ has $`m_1d^{i1}`$ roots in $`W_1`$ and $`m_2d^{i1}`$ roots in $`W_2`$, counting multiplicity. Thus, we may write
$$\varphi ^i(z)a=g_i(z)h_i(z)$$
where $`g_i`$ is monic of degree $`m_1d^{i1}`$ with all its roots in $`W_1`$, and $`h_i`$ is monic of degree $`m_2d^{i1}`$ with all its roots in $`W_2`$. In addition, define $`g_0(z)=h_0(z)=za`$.
We will now use the polynomials $`g_i`$ to compute the bounds given in the Lemma for $`X_1`$; the proof for $`X_2`$ is similar, using $`h_i`$. To simplify notation, write $`X=X_1`$ and $`m=m_1`$.
For any integer $`j0`$, write $`j=c_0+mk`$, and write $`k`$ in base-$`d`$ notation, so that
$$j=c_0+m(c_1+c_2d+c_3d^2+\mathrm{}+c_Md^{M1}),$$
with $`c_0\{0,1,\mathrm{},m1\}`$, and with $`c_i\{0,1,\mathrm{},d1\}`$ for $`i1`$. Define
$$f_j(z)=\underset{i=0}{\overset{M}{}}[g_i(z)]^{c_i}.$$
Clearly, $`f_j`$ is monic of degree $`j`$. Meanwhile, for $`xX`$ and $`i1`$, observe that $`\varphi ^i(x)𝒦_v`$, and therefore $`|\varphi ^i(x)a|r_v`$. On the other hand, all roots of $`h_i`$ lie in $`W_2`$, which is distance $`r_v`$ from $`x`$; therefore, $`|h_i(x)|=r_v^{(dm)d^{i1}}`$. It follows that
$$|g_i(x)|_vr_v^{1(dm)d^{i1}}$$
for all $`i1`$. In addition, since $`XU_0`$, we have $`|g_0(x)|r_v`$. Thus,
$$|f_j(x)|_vr_v^{c_0}\underset{i=1}{\overset{M}{}}r_v^{c_i(1(dm)d^{i1})}=r_v^e,$$
where $`e=e(j,m,d)`$ in the notation of Definition 3.1.
By the same van der Monde determinant argument as in the proof of Lemma 4.1, it follows that if $`N2`$ and $`x_1,\mathrm{},x_NX`$, then
$$\underset{ij}{}|x_ix_j|_vr_v^{E(N,m,d)}.\mathit{}$$
###### Remark 5.2.
In some cases, $`𝒦_v`$ splits into more than two pieces, each much smaller than the $`X_1,X_2`$ of Lemma 5.1. For example, the filled Julia set of the map $`\varphi (z)=z^dc^{d1}z`$ of Example 2.3 breaks naturally into $`d`$ pieces. Adapting the method of the Lemma for each piece, we could ultimately replace the coefficient $`d^22d+2`$ in Theorem 7.1 by $`d`$.
However, as previously noted, most polynomials are not so simple. Indeed, the filled Julia set of $`\varphi (z)=z^daz^{d1}`$ from Example 2.4 splits into only two pieces. (Of course, if we take a higher preimage $`U_n`$ in that example, we get more than two pieces; but because of the large radii, there appears to be no improvement gained by using $`n>1`$.) Even an application of the arguments of Remark 4.3 would result in only a slight decrease in the coefficient of $`N\mathrm{log}_dN`$ in the exponent (cf. Lemma 3.4.b). Unfortunately, a real improvement would require an increase in the size of the (negative) coefficient of $`N^2`$, not the $`N\mathrm{log}_dN`$ term.
## 6. A Partition of the Filled Julia Set: Archimedean Case
The final tool needed for Theorem 7.1 is an archimedean analogue of Lemma 5.1. Roughly the same argument works, but only if the diameter of the filled Julia set $`𝒦`$ is large enough. This phenomenon is familiar to complex dynamicists. For example, given $`\varphi (z)=z^2+c[z]`$, if the diameter of $`𝒦`$ is small, then $`c`$ lies in the Mandelbrot set, in which case $`𝒦`$ is connected. However, once the diameter is large enough, $`c`$ leaves the Mandelbrot set and $`𝒦`$ becomes disconnected. In fact, as the diameter grows, the various pieces of $`𝒦`$ shrink.
We begin with the following preliminary result.
###### Lemma 6.1.
Let $`\varphi [z]`$ be a polynomial of degree $`d2`$ with lead coefficient $`a_d`$. Denote by $`𝒦`$ the filled Julia set of $`\varphi `$ in $``$, and let $`U_0=\overline{D}(a,r^{})`$ be the smallest disk that contains $`𝒦`$. Set $`r=|a_d|^{1/(d1)}r^{}`$, and suppose that
$$r>\{\begin{array}{cc}3,\hfill & \text{if }d=2,\text{ or}\hfill \\ 2+\sqrt{3},\hfill & \text{if }d3.\hfill \end{array}$$
Then $`𝒦`$ is contained in the union of $`d`$ open disks of radius $`|a_d|^{1/(d1)}`$.
###### Proof.
As in the proof of Lemma 4.1, we may assume that $`\varphi `$ is monic. Denote by $`b_1,\mathrm{},b_d`$ the (possibly repeated) roots of $`\varphi (z)=a`$, and let $`\overline{D}(c,s)`$ be the smallest disk containing $`b_1,\mathrm{},b_d`$. (Here, we break our convention and allow $`s=0`$ if $`b_1=\mathrm{}=b_d`$.) Because $`𝒦`$ is not contained in $`D(c,r)`$, there must be some $`y_0𝒦`$ such that $`|y_0c|r`$.
Let $`Y=\overline{D}(c,s)D(y_0,|y_0c|)`$. We claim that $`Y`$ is contained in a disk of radius strictly less than $`s`$ (or that $`Y`$ is empty, if $`s=0`$). Indeed, if $`|y_0c|<s`$, then $`YD(y_0,|y_0c|)`$ trivially. Otherwise, $`|y_0c|s`$, and since the center $`c`$ of the first disk lies on the boundary of the second, the intersection $`Y`$ is contained in a strictly smaller disk. (For example, center the new disk at the midpoint of the two intersection points of the two boundary circles.)
By the minimality of $`s`$, then, not all of $`b_1,\mathrm{},b_d`$ can be in $`Y`$. Thus, there is some $`1id`$ such that $`|y_0b_i||y_0c|r`$. Without loss, assume that $`|y_0b_1|r`$.
For all $`i2`$, we have $`|y_0b_i|rs`$, because $`b_i\overline{D}(c,s)`$. Since $`y_0𝒦`$, we have $`\varphi (y_0)𝒦`$, and therefore $`|\varphi (y_0)a|r`$. If $`rs0`$, then, we have
$$r|\varphi (y_0)a|=\underset{i=1}{\overset{d}{}}|y_0b_i|=|y_0b_1|\underset{i=2}{\overset{d}{}}|y_0b_i|r(rs)^{d1},$$
from which we obtain $`rs1`$. Regardless of the sign of $`rs`$, then, we have $`sr1`$.
Re-index so that $`b_1`$ and $`b_d`$ are distance $`\mathrm{max}\{|b_ib_j|\}`$ apart, and so that for all $`i=1,\mathrm{},d1`$, we have $`|b_{i+1}b_1||b_ib_1|`$. Thus, the $`\{b_i\}`$ are ordered by their distance from $`b_1`$. Moreover, $`|b_db_1|\sqrt{3}s\sqrt{3}(r1)`$; see, for example, , Exercise 6-1.
If $`d=2`$, we can improve this lower bound. In that case, the smallest disk containing $`b_1`$ and $`b_2`$ is the closed disk centered at $`(b_1+b_2)/2`$ of radius $`|b_1b_2|/2`$. That is, $`s=|b_1b_2|/2`$. It follows that $`|b_1b_2|=2s2(r1)`$.
For all degrees $`d2`$, we have $`r>3`$, so that $`s>2`$, and therefore the two disks $`\overline{D}(b_1,1)`$ and $`\overline{D}(b_d,1)`$ are disjoint. Moreover, as $`y`$ ranges through $`[D(b_1,1)D(b_d,1)]`$, the minimum value of $`|yb_1||yb_d|`$ is $`|b_1b_d|1`$, attained at only two points, namely the point on the boundary of each disk closest to the other disk.
Let $`U_1=\varphi ^1(U_0)`$. Since $`𝒦=\varphi ^1(𝒦)U_1`$, it suffices to show that
(6)
$$U_1\underset{i=1}{\overset{d}{}}D(b_i,1).$$
If not, then there is some $`yU_1D(b_i,1)`$. If $`d3`$, then by the above computations, we have $`|yb_1||yb_d|\sqrt{3}(r1)1`$. Since $`\varphi (y)U_0`$, we obtain
$$r|\varphi (y)a|=\underset{i=1}{\overset{d}{}}|yb_i|(\sqrt{3}(r1)1)\underset{i=2}{\overset{d1}{}}|yb_i|(\sqrt{3}(r1)1),$$
contradicting the hypothesis that $`r>2+\sqrt{3}`$. Similarly, if $`d=2`$, then
$$r|\varphi (y)a|=|yb_1||yb_2|>2(r1)1=2r3,$$
contradicting the hypothesis that $`r>3`$, and proving the Lemma. ∎
###### Remark 6.2.
Because $`𝒦_v`$ is compact for archimedean $`v`$, the conclusion of Lemma 6.1 implies that $`𝒦`$ is in fact contained in $`d`$ closed disks of radius strictly less than $`|a_d|_v^{1/(d1)}`$. This fact will be useful in Cases 2 and 3 of the proof of Theorem 7.1.
We are now prepared to present our archimedean version of Lemma 5.1.
###### Lemma 6.3.
Let $`\varphi [z]`$ be a polynomial of degree $`d2`$ with lead coefficient $`a_d`$. Denote by $`𝒦`$ the filled Julia set of $`\varphi `$ in $``$, and let $`r^{}`$ be the radius of the smallest disk $`U_0`$ that contains $`𝒦`$. Set $`r=|a_d|^{1/(d1)}r^{}`$ and
$$C_d=d^{(d2)/(d1)}\mathrm{min}\{1,\frac{1.2}{d1}\}.$$
Suppose that
$$r\{\begin{array}{cc}4\hfill & \text{if }d=2\hfill \\ \frac{\sqrt{3}+2(d1)}{\sqrt{3}(d1)C_d},\hfill & \text{if }d3.\hfill \end{array}$$
Then there are disjoint sets $`X_1,X_2𝒦`$ and positive integers $`m_1,m_2`$ with the properties that $`X_1X_2=𝒦`$, that $`m_1+m_2=d`$, that for $`k=1,2`$, $`\varphi :X_k𝒦_v`$ is $`m_k`$-to-$`1`$, and that for $`k=1,2`$, for any integer $`N2`$, and for any set $`\{x_1,\mathrm{},x_N\}X_k`$ of $`N`$ points in $`X_k`$,
$$\underset{ij}{}|x_ix_j|N^N|a_d|^{N(N1)/(d1)}C_d^{F(N,m_k,d)}\left(C_dr\right)^{E(N,m_k,d)},$$
where $`E(N,m_k,d)`$ and $`F(N,m_k,d)`$ are as in Definition 3.1.
###### Proof.
As in the proof of Lemma 4.1, we may assume that $`\varphi `$ is monic. It is easy to check that $`C_d\mathrm{min}\{1,1.2/(d1)\}`$ (the closest approach for $`d3`$ occurs at $`d=5`$), and that the lower bound $`(\sqrt{3}+2(d1))/(\sqrt{3}(d1)C_d)`$ (respectively, $`4`$) for $`r`$ is greater than $`2+\sqrt{3}`$ (respectively, $`3`$), so that we may invoke Lemma 6.1.
Write $`U_0=\overline{D}(a,r)`$, and define and order $`b_1,\mathrm{},b_d`$ as in the proof of Lemma 6.1, so that $`|b_1b_d|\sqrt{3}(r1)`$ (or $`|b_1b_d|2(r1)`$, if $`d=2`$).
If $`d3`$, observe that for some $`m=1,\mathrm{},d1`$, we have
$$|b_{m+1}b_1||b_mb_1|+2+C_dr.$$
For if not, then
$$\sqrt{3}(r1)|b_db_1|<(d1)\left[2+C_dr\right]=2(d1)+(d1)C_dr,$$
so that $`[\sqrt{3}(d1)C_d]r<\sqrt{3}+2(d1)`$, contradicting the hypotheses.
If $`d=2`$, we have $`|b_2b_1|2r22+r`$, since $`r4`$. Let $`m=1`$ in this case.
Let $`U_1=\varphi ^1(U_0)`$, and set $`W_1=\overline{D}(b_1,|b_mb_1|+1)U_1`$ and $`W_2=U_1W_1`$. Observe that $`\text{dist}(W_1,W_2)C_dr`$. Indeed, if $`y_1W_1`$ and $`y_2W_2`$, then $`y_1\overline{D}(b_i,1)`$ and $`y_2\overline{D}(b_j,1)`$ for some $`1im`$ and some $`m+1jd`$; therefore
$$|y_2y_1||b_jb_1||b_ib_1|2|b_{m+1}b_1||b_mb_1|2C_dr.$$
Since $`W_1`$ contains $`m`$ preimages of $`a`$ and $`W_2`$ contains the other $`dm`$, it follows that $`\varphi `$ maps $`W_1`$ $`m`$-to-$`1`$ onto the connected set $`U_0`$, and it maps $`W_2`$ $`(dm)`$-to-$`1`$ onto $`U_0`$.
Let $`X_1=W_1𝒦`$, $`X_2=W_2𝒦`$, $`m_1=m`$, and $`m_2=dm`$. By the previous paragraph, $`X_1`$ and $`X_2`$ satisfy all of the mapping properties claimed in the Lemma. For any integer $`i0`$, define $`g_i(z)`$ and $`h_i(z)`$ as in the proof of Lemma 5.1. That is, for $`i1`$, write
$$\varphi ^i(z)a=g_i(z)h_i(z),$$
where $`g_i`$ is a monic polynomial of degree $`m_1d^{i1}`$ with all of its roots in $`W_1`$, and $`h_i`$ is a monic polynomial of degree $`m_2d^{i1}`$ with all of its roots in $`W_2`$. For $`i=0`$, define $`g_0(z)=h_0(z)=za`$. We will now compute the bounds given in the Lemma for $`X_1`$; the proof for $`X_2`$ is similar. Write $`X=X_1`$ and $`m=m_1`$.
As in the proof of Lemma 5.1, we may write any integer $`j0`$ as
$$j=c_0+m(c_1+c_2d+c_3d^2+\mathrm{}+c_Md^{M1}),$$
with $`c_0\{0,1,\mathrm{},m1\}`$, and with $`c_i\{0,1,\mathrm{},d1\}`$ for $`i1`$. Similarly, define
$$f_j(z)=\underset{i=0}{\overset{M}{}}[g_i(z)]^{c_i},$$
which is clearly monic of degree $`j`$. As before, for any $`xX`$, we have $`|\varphi ^i(x)a|r`$. Similarly, the roots of $`h_i`$, which all lie in $`W_2`$ (for $`i1`$), are distance at least $`C_dr`$ from $`x`$ (which is worse than the Lemma 5.1 bound of $`r`$). Thus, $`|h_i(x)|(C_dr)^{(dm)d^{i1}}`$, and hence
$$|g_i(x)|r(C_dr)^{(dm)d^{i1}}=C_d^1(C_dr)^{1(dm)d^{i1}}$$
for all $`i1`$. Moreover, since $`xU_0`$, we have $`|g_0(x)|r`$. We obtain
$$|f_j(x)|r^{c_0}\underset{i=1}{\overset{M}{}}C_d^{c_i}(C_dr)^{c_i(1(dm)d^{i1})}=C_d^{(c_0+c_1+\mathrm{}+c_M)}(C_dr)^e,$$
where $`e=e(j,m,d)`$ in the notation of Definition 3.1. The Lemma then follows by the van der Monde determinant argument of the proof of Lemma 4.1. ∎
###### Remark 6.4.
Later, in the proof of Theorem 7.1, we will consider the quantity $`C_dr`$, rather than the radius $`r`$, at the archimedean primes. It is easy to prove that the lower bound for $`r`$ given in Lemma 6.3 is guaranteed to hold provided $`C_dr4+\sqrt{3}`$. (In fact, $`4+\sqrt{3}`$ is the value of $`C_d(\sqrt{3}+2(d1))/(\sqrt{3}(d1)C_d)`$ at $`d=3`$.) For $`d=2`$, we also note the more obvious facts that $`C_2=1`$ and that the corresponding sufficient lower bound for $`C_2r`$ is $`4`$.
###### Remark 6.5.
The bounds in Lemma 6.1 and Lemma 6.3 are not sharp. Besides the fact that most of the comments from Remarks 4.3 and 5.2 apply here, our geometric arguments could also be improved. For example, in the proof of Lemma 6.1, if we considered $`\overline{D}(c,s)D(y_0,t)`$ instead of $`\overline{D}(c,s)D(y_0,|y_0c|)`$, where $`t=\sqrt{|y_0c|^2+s^2}`$, we could show that some $`b_i`$ satisfies $`|y_0b_i|t`$. Related arguments could show that two or more points $`b_i,b_j`$ must make the product $`|y_0b_i||y_0b_j|`$ larger than we proved. Similarly, it should be possible to increase the $`\sqrt{3}`$ factor to something closer to $`2`$ by considering the geometric arrangement of the $`\{b_i\}`$ more delicately.
## 7. The Global Bound
At last, we are prepared to state and prove a precise version of the Main Theorem.
###### Theorem 7.1.
Let $`K`$ be a global field, and let $`\varphi K[z]`$ be a polynomial of degree $`d2`$. Let $`s_{\mathrm{}}0`$ be the number of archimedean primes of $`K`$, and let $`ss_{\mathrm{}}`$ be the number of bad (i.e., not potentially good) primes of $`\varphi `$ in $`M_K`$, including all archimedean primes.
If $`K`$ is a function field, let $`q`$ be the size of the smallest residue field of a prime $`vM_K`$. If $`K`$ is a number field, let $`D=[K\mathrm{:}]`$, and let
$$\sigma =\{\begin{array}{cc}7\hfill & \text{if }d=2,\hfill \\ \frac{233^{(d1)(d2)}}{(d1)(d2)}\hfill & \text{if }d3.\hfill \end{array}$$
Set
$$t=\{\begin{array}{cc}ss_{\mathrm{}}\hfill & \text{if }K\text{ is a number field and }s\sigma D,\hfill \\ s+\frac{D\mathrm{log}d}{2\mathrm{log}2}\hfill & \text{if }K\text{ is a number field and }s>\sigma D,\hfill \\ s\hfill & \text{if }K\text{ is a function field,}\hfill \end{array}$$
and
$$\beta =\{\begin{array}{cc}9\hfill & \text{if }K\text{ is a number field, }s\sigma D,\text{ and }d=2,\hfill \\ \mathrm{max}\{11,2d\}\hfill & \text{if }K\text{ is a number field, }s\sigma D,\text{ and }d3,\hfill \\ 1\hfill & \text{otherwise.}\hfill \end{array}$$
Then $`\varphi `$ has no more than $`M+1`$ $`K`$-rational preperiodic points in $`^1(K)`$, where
$$M=\{\begin{array}{cc}q\hfill & \text{if }K\text{ is a function field and }s=0,\hfill \\ \beta ^D\hfill & \text{if }K\text{ is a number field and }s=s_{\mathrm{}},\hfill \\ \beta ^D(d^22d+2)(t\mathrm{log}_dt+3t)\hfill & \text{if }0<t<d,\hfill \\ \beta ^D(d^22d+2)(t\mathrm{log}_dt+t\mathrm{log}_d\mathrm{log}_dt+3t)\hfill & \text{otherwise.}\hfill \end{array}$$
###### Proof.
For each prime $`vM_K`$, let $`n_v1`$ be the exponent so that the product formula (1) holds for all $`xK^\times `$. Let $`S`$ be the (finite) set of primes of $`K`$ of bad reduction of $`\varphi `$, including all the archimedean primes; that is, $`\mathrm{\#}S=s`$. Let $`a_dK`$ be the lead coefficient of $`\varphi `$. For each prime $`vM_K`$, let $`𝒦_v_v`$ denote the filled Julia set of $`\varphi `$ in $`_v`$, let $`r_v^{}`$ be the radius of the smallest disk in $`_v`$ containing $`𝒦_v`$, and let $`r_v=|a_d|_v^{1/(d1)}r_v^{}`$.
For each non-archimedean prime $`v`$, let $`R_v=r_v^{n_v}`$. For each archimedean prime $`v`$, let $`R_v=(C_dr_v)^{n_v}`$, where $`C_d=d^{(d2)/(d1)}1`$, as in the statement of Lemma 6.3. We consider four cases, some of which overlap with others.
Case 0. The simplest case is that $`K`$ is a function field and $`S=\mathrm{}`$; that is, there are no archimedean primes, and all primes have potentially good reduction. Let $`wM_K`$ be a prime whose residue field has only $`q`$ elements, and suppose that there are $`q+1`$ distinct $`K`$-rational preperiodic points $`\{x_1,\mathrm{},x_{q+1}\}`$ besides the point at $`\mathrm{}`$.
By Lemma 2.5.c, we have $`|x_ix_j|_v|a_d|_v^{1/(d1)}`$ for every $`vM_K`$ and every $`i,j\{1,\mathrm{},n\}`$. Moreover, by the pigeonhole principle, there must be some distinct $`i,j\{1,\mathrm{},n\}`$ such that $`|x_ix_j|_w<|a_d|_w^{1/(d1)}`$. Hence,
$$1=\underset{vM_K}{}\left|x_ix_j\right|_v^{n_v}<\underset{vM_K}{}\left[|a_d|_v^{1/(d1)}\right]^{n_v}=1,$$
which is a contradiction. Thus, there are at most $`q`$ finite $`K`$-rational preperiodic points.
Case 1. Choose $`wM_K`$ such that $`R_wR_v`$ for all $`vM_K`$. (Such $`w`$ exists because $`R_v=1`$ for all but finitely many $`vM_K`$.) In this main case, we suppose that:
* $`R_w>1`$.
* If $`K`$ is a number field, then $`R_w4`$ and $`ss_{\mathrm{}}1`$.
* If $`w`$ is archimedean, then the lower bounds of Lemma 6.3 hold for $`r_w`$.
In particular, we may choose integers $`m_1,m_2`$ and sets $`X_1,X_2𝒦_w`$ for $`\varphi `$ according to Lemma 5.1 (if $`w`$ is non-archimedean) or Lemma 6.3 (if $`w`$ is archimedean).
For each index $`k=1,2`$, set
$$A_k=\frac{dm_k}{m_k(d1)},B_k=1\mathrm{log}_dm_k,\text{and}N_k=M(A_k,B_k,t),$$
where $`M(,,)`$ is as in Lemma 3.5, and where $`t`$ is as in the statement of the Theorem. We claim that there are fewer than $`N_k`$ $`K`$-rational preperiodic points in $`X_k`$.
To prove the claim, fix $`k=1,2`$, and let $`m=m_k`$, $`A=A_k`$, $`B=B_k`$, and $`N=N_k`$. Suppose there are $`N`$ distinct $`K`$-rational preperiodic points $`x_1,\mathrm{},x_N`$ in $`X_k`$. Then by the product formula applied to both $`_{ij}(x_ix_j)`$ and $`a_d`$,
$`1`$ $`={\displaystyle \underset{vM_K}{}}\left|{\displaystyle \underset{ij}{}}(x_ix_j)\right|_v^{n_v}={\displaystyle \underset{vM_K}{}}\left[|a_d|_v^{N(N1)/(d1)}{\displaystyle \underset{ij}{}}|x_ix_j|_v\right]^{n_v}`$
(7) $`{\displaystyle \underset{vS}{}}\left[|a_d|_v^{N(N1)/(d1)}{\displaystyle \underset{ij}{}}|x_ix_j|_v\right]^{n_v},`$
where the inequality is because $`|xy|_v|a_d|_v^{1/(d1)}`$ for all $`vM_KS`$ and $`x,y𝒦_v`$, by Lemma 2.5.c, and because $`x_1,\mathrm{},x_N𝒦_v`$ for every $`vM_K`$.
If $`w`$ is non-archimedean, then by Lemma 4.1 and Lemma 5.1, (7) becomes
$$1N^{DN}r_w^{n_wE(N,m,d)}\underset{vS\{w\}}{}r_v^{n_vE(N,d)}=N^{DN}C_d^{DE(N,d)}R_w^{E(N,m,d)}\underset{vS\{w\}}{}R_v^{E(N,d)},$$
where we set $`D=0`$ if $`K`$ is a function field. (The appearance of $`D`$ in the exponent comes from equation (2).) Because $`R_wR_v`$ and $`E(N,d)0`$, we can replace each $`R_v`$ by $`R_w`$; and because $`R_w,C_d^11`$, we can apply Lemma 3.4.a–b to obtain
(8)
$$1N^{DN}C_d^{(d1)DN\mathrm{log}_dN}R_w^{(d1)N[s\mathrm{log}_dNAN+B]}.$$
Similarly, if $`w`$ is archimedean, then by Lemma 4.1 and Lemma 6.3, (7) becomes
$`1`$ $`N^{DN}C_d^{n_wF(N,M,d)}(C_dr_w)^{n_wE(N,m,d)}{\displaystyle \underset{vS\{w\}}{}}r_v^{n_vE(N,d)}`$
$`=N^{DN}C_d^{(Dn_w)E(N,d)n_wF(N,m,d)}R_w^{E(N,m,d)}{\displaystyle \underset{vS\{w\}}{}}R_v^{E(N,d)}.`$
Replacing each $`R_v`$ by $`R_w`$ as before, and applying Lemma 3.4a–c, we obtain exactly inequality (8) once more.
Meanwhile, we compute
(9)
$$N^{DN}C_d^{(d1)DN\mathrm{log}_dN}=d^{DN\mathrm{log}_dN}d^{(d2)(DN\mathrm{log}_dN)}=d^{(d1)DN\mathrm{log}_dN}.$$
If $`K`$ is a number field, our assumption that $`R_w4`$ means that $`dR_w^{1/\mathrm{log}_d4}`$. Combining (8) and (9), then, we obtain
(10)
$$1R_w^{(d1)N[t\mathrm{log}_dNAN+B]},$$
where $`t=s+D\mathrm{log}d/(2\mathrm{log}2)`$, as in the statement of the Theorem. The same inequality follows for function fields with $`t=s`$, since $`D=0`$ in that case. By our definitions of $`A`$, $`B`$, and $`t`$, the hypotheses of Lemma 3.5 hold. Thus, by that Lemma and our choice of $`N`$, we have $`t\mathrm{log}_dNAN+B<0`$, so that $`1<1`$, which is a contradiction, proving the claim that there are fewer than $`N_k`$ $`K`$-rational preperiodic points in $`X_k`$. (However, since $`N_k`$ need not be an integer, we cannot claim that there are at most $`N_k1`$ such points.)
The total number of finite $`K`$-rational preperiodic points is the number in $`X_1`$ plus the number in $`X_2`$. That is, there are fewer than $`N_1+N_2`$ such points. That upper bound is
(11)
$$N_1+N_2=M(A_1,B_1,t)+M(A_2,B_2,t).$$
From the definition of $`M(A,B,t)`$ in Lemma 3.5, it is each to check that, as $`m_1`$ varies from $`1`$ to $`d1`$, the largest value of $`N_1+N_2`$ in equation (11) is attained at $`m_1=1`$ and $`m_2=d1`$ (or vice versa). In that case, the bound is
$$N_1+N_2=(d^22d+2)\left(t\mathrm{log}_dt+t\mathrm{log}_d(\mathrm{max}\{1,\mathrm{log}_dt\})+3t\right).$$
Adding $`1`$ for the point at $`\mathrm{}`$, we obtain the bound stated in the Theorem, with $`\beta =1`$.
Case 2. Next, suppose that $`K`$ is a number field and $`d=2`$. Write $`S_{\mathrm{}}`$ for the set of archimedean primes of $`M_K`$, and let $`s_{\mathrm{}}=\mathrm{\#}S_{\mathrm{}}`$. We will remove the archimedean primes from the picture by covering the filled Julia set at each such prime $`vS_{\mathrm{}}`$ by at most $`9^{n_v}`$ disks of diameter less than $`|a_d|_v^1`$. To simplify notation, let $`𝒦_v^{}=a_d𝒦_v`$; we wish to cover $`𝒦_v^{}`$ by disks of diameter less than $`1`$.
For any real prime $`vS_{\mathrm{}}`$, the set $`𝒦_v^{}`$ is contained either in a single interval of length $`6`$ or in two intervals of length less than $`2`$, by Lemma 6.1 and Remark 6.2. (In fact, the bound of $`6`$ could be reduced to $`4`$, but we will not need that stronger bound here.) In particular, $`𝒦_v^{}`$ is contained in a union of seven or fewer intervals of length strictly less than $`1`$.
For a complex prime $`vS_{\mathrm{}}`$, the same Lemma implies that $`𝒦_v^{}`$ is contained either in a single disk of radius $`3`$ or in two disks of radius less than $`1`$. Each disk of radius $`1`$ can easily be covered by nine disks of diameter slightly less than $`1`$. Similarly, the disk of radius $`3`$ can be covered by a square of side length $`6`$. That square can then be divided into $`81`$ squares of side length $`2/3`$, each of which fits inside a disk of diameter less than $`1`$.
Scaling back by $`|a_d|_v^1`$, then, we have at each archimedean prime $`vS_{\mathrm{}}`$ at most $`9^{n_v}`$ disks of diameter less than $`|a_d|_v^1`$ which together cover $`𝒦_v`$, as promised. In total, then, we have at most $`9^D`$ choices of one disk for each archimedean prime.
For any such choice $`𝔻=\{D_v:vS_{\mathrm{}}\}`$ of one disk of diameter less than $`|a_d|_v^1`$ for each archimedean prime $`v`$, let $`𝒫_𝔻`$ denote the set of $`K`$-rational preperiodic points $`x`$ for which $`xD_v`$ for every $`vS_{\mathrm{}}`$. We will bound the size of $`𝒫_𝔻`$.
If $`S=S_{\mathrm{}}`$, then each set $`𝒫_𝔻`$ can contain at most one point. Indeed, if there were distinct points $`x,y𝒫_𝔻`$, then
$$1=\underset{vM_K}{}|xy|_v^{n_v}<\underset{vM_K}{}\left[|a_d|_v^1\right]^{n_v}=1,$$
by Lemma 2.5.c, with the strict inequality coming from the fact that the diameter at each archimedean prime is strictly less than $`|a_d|_v^1`$. Since there are $`9^D`$ choices of $`𝔻`$, there are at most $`9^D`$ finite $`K`$-rational preperiodic points.
On the other hand, if $`SS_{\mathrm{}}`$, then choose $`wM_KS_{\mathrm{}}`$ such that $`R_wR_v`$ for all $`vM_KS_{\mathrm{}}`$. By Lemma 2.5.c, $`r_w>1`$, so that we may apply Lemma 5.1 at $`w`$.
Now fix $`𝔻`$ and follow the argument of Case 1, but restricted to $`\{x_i\}𝒫_𝔻`$. At each archimedean prime $`vS_{\mathrm{}}`$ we have $`|x_ix_j|_v|a_d|_v^1`$. Therefore, by Lemmas 4.1, 5.1, and 3.4.a–b, inequality (7) becomes
$$1R_w^{(d1)N[t\mathrm{log}_dNAN+B]},$$
where $`t=ss_{\mathrm{}}`$. Following the rest of the argument of Case 1 (from inequality (10) on), and multiplying by $`9^D`$ (the number of choices $`𝔻`$), we obtain the desired bounds.
Case 3. If $`K`$ is a number field and $`d3`$, we proceed roughly as in Case 2. Again, write $`S_{\mathrm{}}`$ for the set of archimedean primes of $`M_K`$, let $`s_{\mathrm{}}=\mathrm{\#}S_{\mathrm{}}`$, and let $`𝒦_v^{}=\alpha 𝒦_v`$, where $`\alpha ^{d1}=a_d`$. This time, we will cover $`𝒦_v^{}`$ by at most $`\beta ^{n_v}`$ disks of diameter less than $`1`$, where $`\beta =\mathrm{max}\{11,2d\}`$.
For a real prime $`vS_{\mathrm{}}`$, Lemma 6.1 and Remark 6.2 imply that $`𝒦_v^{}`$ is contained either in a single interval of length $`4+2\sqrt{3}`$ or in $`d`$ intervals of length less than $`2`$. In particular, $`𝒦_v^{}`$ is contained in a union of $`\mathrm{max}\{8,2d\}\beta `$ or fewer intervals of length less than $`1`$.
For a complex prime $`vS_{\mathrm{}}`$, the same Lemma implies that $`𝒦_v`$ is contained either in a single disk of radius $`2+\sqrt{3}`$ or in $`d`$ disks of radius less than $`1`$. As before, each disk of radius $`1`$ can be covered by nine disks of diameter less than $`1`$. Similarly, the disk of radius $`2+\sqrt{3}`$ can be covered by a square of side length $`4+2\sqrt{3}`$. That square can be divided into $`121`$ squares of side length $`(4+2\sqrt{3})/11`$, each of which fits inside a disk of diameter less than $`1`$. (In fact, using a hexagonal tiling, one could cover the big disk by $`84`$ disks of diameter less than $`1`$, but the messy proof gives only a minor improvement over $`121`$.) Thus, $`𝒦_v^{}`$ can be covered by a union of $`\mathrm{max}\{121,9d\}\beta ^2`$ disks of diameter less than $`1`$.
The rest of Case 3 then follows Case 2, with $`\beta ^D`$ in place of $`9^D`$. This completes our analysis of the four cases.
Final step. If $`K`$ is a function field, we are done; indeed, by Lemma 2.5.c, Cases 0 and 1 cover all the possibilities.
If $`K`$ is a number field, we will now show that for $`s>\sigma D`$, we are automatically in Case 1. Because $`n_v2`$ for an archimedean prime $`v`$, and by Remark 6.4, we need only show there is some $`wM_K`$ such that $`R_w4^2`$ if $`d=2`$, or such that $`R_w(4+\sqrt{3})^2`$ if $`d3`$.
From basic algebraic number theory, there are at most $`D`$ primes of $`K`$ above any given prime of $``$. Given an integer $`m1`$, let $`p_m`$ denote the $`m^{\text{th}}`$ prime in $``$. (That is, $`p_1=2`$, $`p_2=3`$, $`p_3=5`$, and so on.) Thus, if $`ss_{\mathrm{}}>D(m1)`$, there must be some $`wSS_{\mathrm{}}`$ lying above a prime $`pp_m`$ of $``$. Since $`Ds_{\mathrm{}}`$, we get such a $`w`$ provided $`s>mD`$.
Meanwhile, by Lemma 2.8, given $`wSS_{\mathrm{}}`$ lying over a prime $`p`$ of $``$, we have
$$R_w|\pi _w|_w^{n_w}p\text{if }d=2,$$
or
$$R_w^{(d1)(d2)}|\pi _w|_w^{n_w}p\text{if }d3,$$
where $`\pi _w`$ is a uniformizer at $`w`$. (The Lemma applies because if $`𝒦_wK=\mathrm{}`$, then there are no finite $`K`$-rational preperiodic points at all, and the conclusion of the Theorem is trivial.) For $`d=2`$, then, the condition $`R_w16`$ is guaranteed provided $`s7D+1`$, since $`17`$ is the seventh prime of $``$. Thus, $`s>7D=\sigma D`$ suffices for $`d=2`$.
For $`d3`$, the elementary estimate in Theorem 4.7 of says that $`p_m>(1/6)m\mathrm{log}m`$ for any integer $`m1`$. It is easy to check that $`m=\sigma `$ satisfies $`m\mathrm{log}m6(4+\sqrt{3})^{2(d1)(d2)}`$, where $`\sigma =233^{(d1)(d2)}/[(d1)(d2)]`$ as in the statement of the Theorem. (The $`33`$ appears because it is the smallest integer larger than $`(4+\sqrt{3})^2`$.) Thus, $`s>\sigma D`$ implies $`R_w(4+\sqrt{3})^2`$, once again forcing Case 1. ∎
###### Remark 7.2.
If, for a given polynomial $`\varphi `$, we know that we are in Case 1 (say, by inspection of the filled Julia set at one prime), then we can set $`\beta =1`$ in the statement of the Theorem, even if $`s\sigma D`$. In particular, for a fixed function $`\varphi `$, the conditions of Case 1 are preserved if one passes to a finite extension of $`K`$. Thus, one would not have to worry about the growth of $`s`$ relative to $`\sigma D`$ as one traveled up a tower of number fields, even though one cannot expect $`s`$ to increase as fast as $`\sigma D`$ in general.
###### Remark 7.3.
Our covering methods in Cases 2 and 3 are rather crude, and it should be possible to cover the filled Julia sets $`𝒦_v`$ in these cases far more efficiently. For example, our use of disks of diameter $`1`$ was rather simplistic. Instead, one could cover $`𝒦_v`$ by larger sets $`Y`$ for which $`_{1i,jL}|y_iy_j|_v1`$ for some fixed small integer $`L`$. Such modifications could lead to a substantial reduction in the coefficient $`\beta ^D`$ in the final bound.
Even without any extra work, the coefficient can be improved in special cases. For example, if $`K`$ is a totally real number field, then the cutoff $`\sigma `$, which determines when $`\beta `$ drops to $`1`$, would be much smaller, since we would only need $`R_w4`$ (if $`d=2`$), or $`R_w4+\sqrt{3}`$ (if $`d3`$) rather than $`4^2`$ or $`(4+\sqrt{3})^2`$. Moreover, if $`K`$ is totally real and $`d=2`$, then each archimedean filled Julia set is contained in a union of four intervals of length $`1`$. (See, for example, Lemma 6.4 and Proposition 6.6 of .) Thus, the coefficient of $`9^D`$ that appears in Theorem 7.1 could be replaced by $`4^D`$, with one exception.
The one exception is if all non-archimedean primes have good reduction and the archimedean filled Julia set is an interval of length $`4`$. This occurs for the Chebyshev polynomial $`\varphi (z)=z^22`$, which has filled Julia set $`[2,2]`$. In this special case, after removing the points $`\mathrm{}`$ and $`2`$, the rest of the preperiodic points can be covered by four half-open intervals of length $`1`$ at each archimedean prime. Since there are no non-archimedean bad primes, we obtain a bound of $`2+4^D`$ for the total number of preperiodic points in $`^1(K)`$.
###### Remark 7.4.
Another approach to finding a cutoff $`\sigma `$ which forces Case 1 would be to consider the set $`TS`$ consisting of non-archimedean bad primes $`v`$ at which there are actually $`K`$-rational preperiodic points $`x,y`$ for which $`|xy|_v>1`$. For such primes, the exponent of $`1/[(d1)(d2)]`$ in Lemma 2.8 could be improved to $`1/(d1)`$. Unfortunately, there may not be very many such primes. As a result, although the exponent of $`(d1)(d2)`$ in the definition of $`\sigma `$ could be improved to $`(d1)`$, it would come at the expense of introducing a extra factor like $`\beta ^D`$ into the formula for $`\sigma `$.
###### Remark 7.5.
For large degrees $`d`$, one can obtain slightly smaller bounds by using more than one big bad prime $`w`$. There is, of course, a trade-off. While using $`\mathrm{}2`$ big primes $`w`$ ultimately increases the coefficient $`A`$ of $`N`$ in the exponent of (8), it also increases the number of pieces $`\{X_k\}`$ from $`2`$ to $`2^{\mathrm{}}`$. It appears that the optimal number of such primes to use is $`\mathrm{}2\mathrm{log}_2(d1)`$. The improved bound for the number of rational preperiodic points would be roughly the old bound divided by $`2\mathrm{log}_2(d1)`$, for large $`d`$. However, the proof would be vastly more complicated, especially dealing with the archimedean primes. The slight improvement seems not to be worth the increased difficulty, especially given that the resulting bound would still be very far from the conjectured uniform bound.
We close by presenting a slight strengthening of Theorem 7.1 in the simplest case.
###### Example 7.6.
Let $`K=`$ (so that $`D=s_{\mathrm{}}=1`$, and $`n_v=1`$ for all $`vM_{}`$) and $`d=2`$. That is, we wish to bound the number of rational preperiodic points of a quadratic polynomial $`\varphi [z]`$. It is of course well known that any such polynomial is conjugate over $``$ to one of the form $`\varphi _c(z)=z^2+c`$, with $`c`$.
Let us suppose that $`\varphi _c`$ has at least one preperiodic point in $``$. This supposition implies that $`c=j/m^2`$ for some relatively prime integers $`j,m`$, and that $`\mathrm{}<c1/4`$; see, for example, Proposition 6.7 of . (One can also easily establish that $`j`$ must satisfy one of approximately $`2^s`$ congruences modulo $`m`$, but we do not need that here.) For non-archimedean primes $`v`$ of $``$, we have $`R_v=|m|_v^1`$ if $`v`$ is odd, and $`R_2=\mathrm{max}\{|m/2|_2^1,1\}`$. (Note that if $`4m`$, then $`\varphi _c`$ has good reduction at $`v=2`$, after a change of coordinates.) In addition, for $`c<0`$, we have $`R_{\mathrm{}}=(1+\sqrt{14c})/2`$.
We will be in Case 1 of the proof of Theorem 7.1 provided there is some prime $`v`$ with $`R_v4`$. Still assuming that there is at least one preperiodic point in $``$, Lemma 2.8 says that such a prime must exist unless the only bad primes are $`\mathrm{}`$, $`2`$, and $`3`$, and $`R_2,R_3,R_{\mathrm{}}<4`$. By our characterization of $`R_v`$ above, this means that the denominator $`m`$ is a divisor of $`12`$, and that $`12<c1/4`$. There are only finitely many rational numbers of the form $`c=j/144`$ between $`12`$ and $`1/4`$, and a simple computer search shows none of the corresponding polynomials $`\varphi _c`$ has more than eight preperiodic points in $``$. (For five such values of $`c`$, namely $`21/16`$, $`29/16`$, $`91/36`$, $`133/144`$, and $`1333/144`$, there are exactly eight preperiodic points in $``$. Incidentally, there are infinitely many values $`c`$ for which $`\varphi _c`$ has at least eight preperiodic point in $``$, by Theorem 2 of .)
For all other $`c`$, we are in Case 1, which means $`t=s+D\mathrm{log}d/(2\mathrm{log}2)=s+1/2`$, and $`\beta =1`$. If $`s=1`$, then only the archimedean prime is bad, and in light of Remark 7.3, there are at most five preperiodic points in $``$; in fact, there are at most four for $`s=1`$ and $`c2`$. The only remaining possibility is that $`s2`$ and $`\beta =1`$, in which case the number of preperiodic points in $``$ is at most
$$(2s+1)\left[\mathrm{log}_2(2s+1)+\mathrm{log}_2(\mathrm{log}_2(2s+1)1)+2\right].$$
Since this bound is greater than eight even for $`s=2`$, it holds even without making the exceptions from the previous paragraph. |
warning/0506/math0506029.html | ar5iv | text | # Dynamic Integration of Time- and State-domain Methods for Volatility Estimation
## 1 Introduction
In forecasting a future event or making an investment decision, two pieces of useful information are frequently consulted. Based on the recent history, one uses a form of local average, such as the moving average in the time-domain, to forecast a future event. This approach uses the continuity of a function and ignores completely the information in the remote history, which is related to current through stationarity. On the other hand, one can forecast a future event based on state-domain modeling such as the ARMA, TAR, ARCH models or nonparametric models (see Tong, 1990; Fan & Yao, 2003 for details). For example, to forecast the volatility of the yields of a bond with the current rate 6.47%, one computes the standard deviation based on the historical information with yields around 6.47%. This approach relies on the stationarity and depends completely on historical data. But, it ignores the importance of the recent data. The question of how to combine the estimators from both the time-domain and the state-domain poses an interesting challenge to statisticians.
To elucidate our idea, consider the weekly data on the yields of 3-month treasury bills presented in Figure 1. Suppose that the current time is January 04, 1991 and interest rate is 6.47% on that day, corresponding to the time index $`t=1930`$. One may estimate the volatility based on the weighted squared differences in the past 52 weeks (1 year), say. This corresponds to the time-domain smoothing, using a small vertical stretch of data in Figure 1(a). Figure 1(b) computes the squared differences of the past year’s data and depicts its associated exponential weights. The estimated volatility (conditional variance) is indicated by the dashed horizontal bar. Let the resulting estimator be $`\widehat{\sigma }_{t,\text{time}}^2`$. On the other hand, in financial activities, we do consult historical information in making better decisions. The current interest rate is 6.47%. One may examine the volatility of the yields when the interest rates are around 6.47%, say, $`6.47\%\pm .25\%`$. This corresponds to using the part of data indicated by the horizontal bar. Figure 1(c) plots the squared differences $`X_tX_{t1}`$ against $`X_{t1}`$ with $`X_{t1}`$ restricted to the interval $`6.47\%\pm .25\%`$. Applying the local kernel weight to the squared differences results in a state-domain estimator $`\widehat{\sigma }_{t,\text{state}}^2`$, indicated by the horizontal bar in Figure 1(c). Clearly, as shown in Figure 1(a), except in the 3-week period right before January 4, 1991 (which can be excluded in the state domain fitting), the last period with interest rate around $`6.47\%\pm .25\%`$ is the period from May 15, 1988 and July 22, 1988. Hence, the time and state-domain estimators use two nearly independent components of the time series, as they are 136-week apart in time. See the horizontal and vertical bars of Figure 1(a). These two kinds of estimators have been used in the literature for forecasting volatility. The former is prominently featured in the RiskMetrics of J.P. Morgan, and the latter has been used in nonparametric regression (see Tong, 1995; Fan & Yao, 2003 and references therein). The question arises how to integrate them.
An integrated estimator is to introduce a dynamic weighting scheme $`0w_t1`$ to combine the two nearly independent estimators. Define the resulting integrated estimators as
$$\widehat{\sigma }_t^2=w_t\widehat{\sigma }_{t,\text{time}}^2+(1w_t)\widehat{\sigma }_{t,\text{state}}^2.$$
The question is how to choose the dynamic weight $`w_t`$ to optimize the performance. A reasonable approach is to minimize the variance of the combined estimator, leading to the dynamic optimal weights
$$w_t=\frac{\text{Var}(\widehat{\sigma }_{t,\text{state}}^2)}{\text{Var}(\widehat{\sigma }_{t,\text{time}}^2)+\text{Var}(\widehat{\sigma }_{t,\text{state}}^2)},$$
(1)
since the two piece of estimators are nearly independent. The unknown variances in (1) can easily be estimated in Section 3. Another approach is the Bayesian approach, which regards the historical information as the prior. We will explore this idea in Section 4. The proposed method is also applicable to other estimation problems in time series such as forecasting the mean function and the volatility matrix of multivariate time series.
To appreciate the intuition behind our approach, let us consider the diffusion process
$$dr_t=\mu (r_t)dt+\sigma (r_t)dW_t,$$
(2)
where $`W_t`$ is a Wiener process. This diffusion process is frequently used to model asset price and the yields of bonds, which are fundamental to fixed income securities, financial markets, consumer spending, corporate earnings, asset pricing and inflation. The family of models include famous ones such as the Vasicek (1977) model, the CIR model (Cox, et al. 1985) and the CKLS model (Chan, et al. 1992). Suppose that at time $`t`$ we have a historic data $`\{r_{t_i}\}_{i=0}^N`$ from the process (2) with a sampling interval $`\mathrm{\Delta }`$. Our aim is to estimate the volatility $`\sigma _t^2\sigma ^2(r_t).`$ Let $`Y_i=\mathrm{\Delta }^{1/2}(r_{t_{i+1}}r_{t_i})`$. Then for the model (2), the Euler approximation scheme is
$$Y_i\mu (r_{t_i})\mathrm{\Delta }^{1/2}+\sigma (r_{t_i})\epsilon _i,$$
(3)
where $`\epsilon _i_{i.i.d.}N(0,1)`$ for $`i=0,\mathrm{},N1`$. Fan & Zhang (2003) studied the impact of the order of difference on statistical estimation. They found that while higher order can possibly reduce approximation errors, it increases variances of data substantially. They recommended the Euler scheme (3) for most practical situations. The time-domain smoothing relies on the smoothness of $`\sigma (r_{t_i})`$ as a function of time $`t_i`$. This leads to the exponential smoothing estimator in Section 2.1. On the other hand, the state-domain smoothing relies on structural invariability implied by the stationarity: the conditional variance of $`Y_i`$ given $`r_{t_i}`$ remains the same even for the data in the history. In other words, historical data also furnish the information about $`\sigma ()`$ at the current time. Combining these two nearly independent estimators leads to a better estimator.
In this paper, we focus on the estimation of volatility of a portfolio to illustrate how to deal with the problem of dynamic integration. Asymptotic normality of the proposed estimator is established and extensive simulations are conducted, which theoretically and empirically demonstrate the dominated performance of the integrated estimation.
## 2 Estimation of Volatility
The volatility estimation is an important issue of modern financial analysis since it pervades almost every facet of this field. It is a measure of risk of a portfolio and is related to the Value-at-Risk (VaR), asset pricing, portfolio allocation, capital requirement and risk adjusted returns, among others. There is a large literature on estimating the volatility based on time-domain and state-domain smoothing. For an overview, see the recent book by Fan & Yao (2003).
### 2.1 Time-domain estimator
A popular version of time-domain estimator of the volatility is the moving average estimator:
$$\widehat{\sigma }_{MA,t}^2=n^1\underset{i=tn}{\overset{t1}{}}Y_i^2,$$
(4)
where $`n`$ is the size of the moving window. This estimator ignores the drift component, which contributes to the variance in the order of $`O(\mathrm{\Delta })`$ instead of $`O(\mathrm{\Delta }^{1/2})`$ (see Stanton, 1997 and Fan & Zhang, 2003), and utilizes local $`n`$ data points. An extension of the moving average estimator is the exponential smoothing estimation of the volatility given by
$$\widehat{\sigma }_{ES,t}^2=(1\lambda )Y_{t1}^2+\lambda \widehat{\sigma }_{ES,t1}^2=(1\lambda )\{Y_{t1}^2+\lambda Y_{t2}^2+\lambda ^2Y_{t3}^2+\mathrm{}\},$$
(5)
where $`\lambda `$ is a smoothing parameter that controls the size of the local neighborhood. The RiskMetrics of J.P. Morgan (1996), which is used for measuring the risks, called Value at Risk (VaR), of financial assets, recommends $`\lambda =0.94`$ and $`\lambda =0.97`$ respectively for calculating VaR of the daily and monthly returns.
The exponential smoothing estimator in (5) is a weighted sum of the squared returns prior to time $`t`$. Since the weight decays exponentially, it essentially uses recent data. A slightly modified version that explicitly uses only $`n`$ data points before time $`t`$ is
$$\widehat{\sigma }_{ES,t}^2=\frac{1\lambda }{1\lambda ^n}\underset{i=1}{\overset{n}{}}Y_{ti}^2\lambda ^{i1}.$$
(6)
When $`\lambda =1`$, it becomes the moving average estimator (1). With slight abuse of notation, we will also denote the estimator for $`\sigma ^2(r_t)`$ as $`\widehat{\sigma }_{ES,t}^2`$.
All of the time domain smoothing is based on the assumption that the returns $`Y_{t1},Y_{t2},`$ $`\mathrm{},Y_{tn}`$ have approximately the same volatility. In other words, $`\sigma (r_t)`$ in (1) is continuous in time $`t`$. The following proposition gives the condition under which this holds.
###### Proposition 1
Under Conditions (A1) and (A2) in the Appendix, we have
$$|\sigma ^2(r_s)\sigma ^2(r_u)|K|su|^{(p1)/(2p)},$$
for any $`s,u[t\eta ,t]`$, where the coefficient $`K`$ satisfies $`E[K^{2(p+\delta )}]<\mathrm{}`$ and $`\eta `$ is a positive constant.
With the above Hölder continuity, we can establish the asymptotic normality of the time-domain estimator.
###### Theorem 1
Suppose that $`\sigma _t^2>0`$. Under conditions (A1) and (A2), if $`n+\mathrm{}`$ and $`n\mathrm{\Delta }0`$, then
$$\widehat{\sigma }_{ES,t}^2\sigma _t^20,\text{a.e.}$$
Moreover, if the limit $`c=lim_n\mathrm{}n(1\lambda )`$ exists and $`n\mathrm{\Delta }^{(p1)/(2p1)}0`$,
$$\sqrt{n}[\widehat{\sigma }_{ES,t}^2\sigma _t^2]/s_{1,t}\stackrel{𝒟}{}𝒩(0,1),$$
where $`s_{1,t}^2=c\sigma _t^4\frac{e^c+1}{e^c1}.`$
Theorem 1 has very interesting implications. Even though the data in the local time-window is highly correlated (indeed, the correlation tending to one), we can compute the variance as if the data were independent. Indeed, if the data in (6) were independent and locally homogeneous, we have
$`\text{Var}(\widehat{\sigma }_{ES,t}^2)`$ $``$ $`{\displaystyle \frac{(1\lambda )^2}{(1\lambda ^n)^2}}2\sigma _t^4{\displaystyle \underset{i=1}{\overset{n}{}}}\lambda ^{2(i1)}`$
$`=`$ $`{\displaystyle \frac{2\sigma _t^4(1\lambda )(1+\lambda ^n)}{(1+\lambda )(1\lambda ^n)}}{\displaystyle \frac{1}{n}}s_{1,t}^2.`$
This is indeed the asymptotic variance given in Theorem 1.
### 2.2 Estimation in state-domain
To obtain the nonparametric estimation of the functions $`f(x)=\mathrm{\Delta }^{1/2}\mu (x)`$ and $`\sigma ^2(x)`$ in (3), we use the local linear smoother studied in Ruppert et al. (1997) and Fan & Yao (1998). The local linear technique is chosen for its several nice properties, such as the asymptotic minimax efficiency and the design adaptation. Further, it automatically corrects edge effects and facilitates the bandwidth selection (Fan & Yao, 2003).
To facilitate the theoretical argument in Section 3, we exclude the $`n`$ data points used in the time-domain fitting. Thus, the historical data at time $`t`$ are $`\{(r_{t_i},Y_i),i=0,\mathrm{},Nn1\}`$. Let $`\widehat{f}(x)=\widehat{\alpha }_1`$ be the local linear estimator that solves the following weighted least-squares problem:
$$(\widehat{\alpha }_1,\widehat{\alpha }_2)=\mathrm{arg}\underset{\alpha _1,\alpha _2}{\mathrm{min}}\underset{i=0}{\overset{Nn1}{}}[Y_i\alpha _1\alpha _2(r_{t_i}x)]^2K_{h_1}(r_{t_i}x),$$
where $`K()`$ is a kernel function and $`h_1>0`$ is a bandwidth. Denote the squared residuals by $`\widehat{R}_i=\{Y_i\widehat{f}(r_{t_i})\}^2`$. Then the local linear estimator of $`\sigma ^2(x)`$ is $`\widehat{\sigma }_S^2(x)=\widehat{\beta }_0`$ given by
$$(\widehat{\beta }_0,\widehat{\beta }_1)=\mathrm{arg}\underset{\alpha ,\beta }{\mathrm{min}}\underset{i=0}{\overset{Nn1}{}}\{\widehat{R}_i\beta _0\beta _1(r_{t_i}x)\}^2W_h(r_{t_i}x)$$
(7)
with kernel function $`W`$ and bandwidth $`h`$. Fan & Yao (1998) gives strategies of bandwidth selection. It was shown in Stanton (1997) and Fan & Zhang (2003) that $`Y_i^2`$ instead of $`\widehat{R}_i`$ in (7) can also be used for the estimation of $`\sigma ^2(x)`$.
The asymptotic bias and variance of $`\widehat{\sigma }_S^2(x)`$ are given by Fan & Zhang (2003, theorem 4). Set $`\nu _j=u^jW^2(u)𝑑u`$ for $`j=0,1,2`$. Let $`p()`$ the invariant density function of the Markov process $`\{r_s\}`$ from (1). Then, we have
###### Theorem 2
Let $`x`$ be in the interior of the support of $`p()`$. Suppose that the second derivatives $`\mu ()`$ and $`\sigma ^2()`$ exist in a neighborhood of $`x`$. Under conditions (A3)-(A7), we have
$`\sqrt{(Nn)h}[\widehat{\sigma }_S^2(x)\sigma ^2(x)]/s_2(x)\stackrel{𝒟}{}𝒩(0,1),`$
where $`s_2^2(x)=2\nu _0\sigma ^4(x)/p(x).`$
## 3 Dynamic Integration of time and state domain estimators
In this section, we first show how the optimal dynamic weights in (1) can be estimated and then prove that the time-domain and state-domain estimator are indeed asymptotically independent.
### 3.1 Estimation of dynamic weights
For the exponential smoothing estimator in (6), we can apply the asymptotic formula given in Theorem 1 to get an estimate of its asymptotic variance. However, since the estimator is a weighted average of $`Y_{ti}^2`$, we can obtain its variance directly by assuming $`Y_{tj}N(0,\sigma _t^2)`$ for small $`j`$. Indeed, with the above local homogeneous model, we have
$`\text{Var}(\widehat{\sigma }_{ES,t}^2)`$ $``$ $`{\displaystyle \frac{(1\lambda )^2}{(1\lambda ^n)^2}}2\sigma _t^4{\displaystyle \underset{i=1}{\overset{n}{}}}{\displaystyle \underset{j=1}{\overset{n}{}}}\lambda ^{i+j2}\rho (|ij|)`$ (8)
$`=`$ $`{\displaystyle \frac{2(1\lambda )^2\sigma _t^4}{(1\lambda ^n)^2}}\{1+2{\displaystyle \underset{k=1}{\overset{n1}{}}}\rho (k)\lambda ^k(1\lambda ^{2(nk)})/(1\lambda ^2)\},`$
where $`\rho (j)=\text{Cor}(Y_t^2,Y_{tj}^2)`$ is the autocorrelation of the series $`\{Y_{tj}^2\}`$. The autocorrelation can be estimated from the data in history. Note that due to the locality of the exponential smoothing, only $`\rho (j)`$’s with the first 30 lags, say, contribute to the variance calculation.
We now turn to estimate the variance of $`\widehat{\sigma }_{S,t}^2=\widehat{\sigma }_S^2(r_t)`$. Details can be found in Fan & Yao (1998) and §6.2 of Fan & Yao (2003). Let
$$V_j(x)=\underset{i=1}{\overset{t1}{}}(r_{t_i}x)^jW\left(\frac{r_{t_i}x}{h_1}\right)$$
and
$$\xi _i(x)=W\left(\frac{r_{t_i}x}{h_1}\right)\{V_2(x)(r_{t_i}x)V_1(x)\}/\{V_0(x)V_2(x)V_1(x)^2\}.$$
Then the local linear estimator can be expressed as
$$\widehat{\sigma }_S^2(x)=\underset{i=1}{\overset{t1}{}}\xi _i(x)\widehat{R}_i$$
and its variance can be approximated as
$$\text{Var}(\widehat{\sigma }_S^2(x))\text{Var}\{(Y_1f(x))^2|r_{t_1}=x\}\underset{i=1}{\overset{t1}{}}\xi _i^2(x).$$
(9)
See also Figure 1 and the discussions at the end of §2.1. Again, for simplicity, we assume that $`\text{Var}(\widehat{R}_i|r_{t_i}=x)2\sigma ^4(x)`$, which holds if $`\epsilon _tN(0,1)`$.
Combining (1), (8) and (9), we propose to combine the time-domain and the state-domain estimator with the dynamic weight
$$\widehat{w}_t=\frac{\widehat{\sigma }_{S,t}^4_{i=1}^{t1}\xi _i^2(r_t)}{\widehat{\sigma }_{S,t}^4_{i=1}^{t1}\xi _i^2(r_t)+c_t\widehat{\sigma }_{ES,t}^4},$$
(10)
where $`c_t=\frac{(1\lambda )^2}{(1\lambda ^n)^2}\{1+2_{k=1}^{n1}\rho (k)\lambda ^k(1\lambda ^{2(nk)})/(1\lambda ^2)\}`$ \[see (8)\]. This is obtained by substituting (8) and (9) into (1). For practical implementation, we truncate the series $`\{\rho (i)\}_{i=1}^{t1}`$ in the summation as $`\{\rho (i)\}_{i=1}^{30}`$. This results in the dynamically integrated estimator
$$\widehat{\sigma }_{I,t}^2=\widehat{w}_t\widehat{\sigma }_{ES,t}^2+(1\widehat{w}_t)\widehat{\sigma }_{S,t}^2,$$
(11)
where $`\widehat{\sigma }_{S,t}^2=\widehat{\sigma }_S^2(r_t)`$. The function $`\widehat{\sigma }_S^2()`$ depends on the time $`t`$ and we need to update this function as time evolves. Fortunately, we need only to know the function at the point $`r_t`$. This reduces significantly the computational cost. The computational cost can be reduced further, if we update the estimated function $`\widehat{\sigma }_{S,t}^2`$ at a prescribed time schedule (e.g. once every two months for weekly data).
Finally, we would like to note that in the choice of weight, only the variance of the estimated volatility is considered, rather than the mean square error. This is mainly to facilitate the dynamically weighted procedure. Since the smoothing parameters in $`\widehat{\sigma }_{ES,t}^2`$ and $`\widehat{\sigma }_S^2(x)`$ have been tuned to optimize their performance separately, their biases and variances trade-off have been considered. Hence, controlling the variance of the integrated estimator $`\widehat{\sigma }_{I,t}^2`$ has also controlled, to some extent, the bias of the estimator. Our method focuses only on the estimation of volatility, but the method can be adapted to other estimation problems, such as the value at risk studied in Duffie & Pan (1997) and the drift estimation for diffusion considered in Spokoiny (2000) and volatility matrix for multivariate time series. Further study along this topic is beyond the scope of the current investigation.
### 3.2 Sampling properties
The fundamental component to the choice of dynamic weights is the asymptotic independent between the time and state-domain estimator. By ignoring the drift term (see Stanton, 1997; Fan & Zhang 2003), both the estimators $`\widehat{\sigma }_{ES,t}^2`$ and $`\widehat{\sigma }_{S,t}^2`$ are linear in $`\{Y_i^2\}`$. The following theorem shows that the time-domain and state-domain estimators are indeed asymptotically independent. To facilitate the notation, we present the result at the current time $`t_N`$.
###### Theorem 3
Let $`s_{2,t_N}=s_2(r_{t_N}).`$ Under the conditions of Theorems 1 and 2, if the condition (A2) holds at point $`t_N`$, we have
* asymptotic independence:
$$[\sqrt{n}(\widehat{\sigma }_{ES,t_N}^2\sigma _{t_N}^2)/s_{1,t_N},\sqrt{(Nn)h}(\widehat{\sigma }_{S,t_N}^2\sigma _{t_N}^2)/s_{2,t_N}]^T\stackrel{𝒟}{}𝒩(0,I_2).$$
* asymptotic normality of $`\widehat{\sigma }_{I,t_N}^2`$: if the limit $`d=lim_N\mathrm{}n/[(Nn)h]`$ exists, then
$$\sqrt{(Nn)h/\omega }[\widehat{\sigma }_{I,t_N}^2\sigma _{t_N}^2)]\stackrel{𝒟}{}𝒩(0,1),$$
where $`\omega =w_{t_N}^2s_{1,t_N}^2/d+(1w_{t_N})^2s_{2,t_N}^2`$.
From Theorem 3, based on the optimal weight the asymptotic relative efficiencies of $`\widehat{\sigma }_{I,t_N}^2`$ with respect to $`\widehat{\sigma }_{S,t_N}^2`$ and $`\widehat{\sigma }_{ES,t_N}^2`$ are respectively
$$\text{eff}(\widehat{\sigma }_{I,t_N}^2,\widehat{\sigma }_{S,t_N}^2)=1+ds_{2,t_N}^2/s_{1,t_N}^2,\text{and}\text{eff}(\widehat{\sigma }_{I,t_N}^2,\widehat{\sigma }_{ES,t_N}^2)=1+s_{1,t_N}^2/(ds_{2,t_N}^2),$$
which are greater than one. This demonstrates that the integrated estimator $`\widehat{\sigma }_{I,t_N}^2`$ is more efficient than the time domain and the state domain estimators.
## 4 Bayesian integration of volatility estiamtes
Another possible approach is to consider the historical information as the prior and to incorporate them in the estimation of volatility by the Bayesian framework. We now explore such an approach.
### 4.1 Bayesian estimation of volatility
The Bayesian approach is to regard the recent data $`Y_{tn},\mathrm{},Y_{t1}`$ as an independent sample from $`N(0,\sigma ^2)`$ \[see (3)\] and to regard the historical information being summarized in a prior. To incorporate historical information, we assume that the variance $`\sigma ^2`$ follows an Inverse Gamma distribution with parameters $`a`$ and $`b`$, which has the density function
$$f(\sigma ^2)=b^a\mathrm{\Gamma }^1(a)\{\sigma ^2\}^{(a+1)}\text{exp}(b/\sigma ^2).$$
Denote by $`\sigma ^2IG(a,b)`$. It is a well-known fact that
$$\text{E}(\sigma ^2)=\frac{b}{(a1)},\text{Var}(\sigma ^2)=\frac{b^2}{(a1)^2(a2)},\text{mode}(\sigma ^2)=\frac{b}{(a+1)}.$$
(12)
The hyperparameters $`a`$ and $`b`$ will be estimated from historical data such as the state-domain estimators.
It can easily be shown that the posterior density of $`\sigma ^2`$ given $`\text{Y}=(Y_{tn},\mathrm{},Y_{t1})`$ is IG$`(a^{},b^{})`$, where
$$a^{}=a+\frac{n}{2},b^{}=\frac{1}{2}\underset{i=1}{\overset{n}{}}Y_{ti}^2+b.$$
From (12), the Bayesian mean of $`\sigma ^2`$ is
$$\widehat{\sigma }^2=\frac{b^{}}{(a^{}1)}=\underset{i=1}{\overset{n}{}}(Y_{ti}^2+2b)/(2(a1)+n).$$
This Bayesian estimator can easily be written as
$$\widehat{\sigma }_B^2=\frac{n}{n+2(a1)}\widehat{\sigma }_{MA,t}^2+\frac{2(a1)}{n+2(a1)}\widehat{\sigma }_P^2,$$
(13)
where $`\widehat{\sigma }_{MA,t}^2`$ is the moving average estimator given by (4) and $`\widehat{\sigma }_P^2=b/(a1)`$ is the prior mean, which will be determined from the historical data. This combines the estimate based on the data and prior knowledge.
The Bayesian estimator (14) utilizes the local average of $`n`$ data points. To incorporate the exponential smoothing estimator (5), we regard it as the local average of
$$n^{}=\underset{i=1}{\overset{n}{}}\lambda ^{i1}=\frac{1\lambda ^n}{1\lambda }$$
(14)
data points. This leads to the following integrated estimator
$`\widehat{\sigma }_{B,t}^2`$ $`=`$ $`{\displaystyle \frac{n^{}}{n^{}+2(a1)}}\widehat{\sigma }_{ES,t}^2+{\displaystyle \frac{2(a1)}{2(a1)+n^{}}}\widehat{\sigma }_P^2`$ (15)
$`=`$ $`{\displaystyle \frac{1\lambda ^n}{1\lambda ^n+2(a1)(1\lambda )}}\widehat{\sigma }_{ES,t}^2+{\displaystyle \frac{2(a1)(1\lambda )}{1\lambda ^n+2(a1)(1\lambda )}}\widehat{\sigma }_P^2.`$
In particular, when $`\lambda =1`$, the estimator (15) reduces to (13).
### 4.2 Estimation of Prior Parameters
A reasonable source for obtaining the prior information in (15) is based on the historical data up to time $`t`$. Hence, the hyper-parameters $`a`$ and $`b`$ should depend on $`t`$ and can be used to match with the historical information. Using the approximation model (3), we have
$$E[(Y_t\widehat{f}(r_t))^2r_t]\sigma ^2(r_t)\text{Var}[(Y_t\widehat{f}(r_t))^2r_t]2\sigma ^4(r_t).$$
These can be estimated from the historical data up to time $`t`$, namely, the state-domain estimator $`\widehat{\sigma }_S^2(r_t)`$. Since we have assumed that prior distribution for $`\sigma _t^2`$ is IG($`a_t,b_t)`$, then by the method of moments, we would set
$`E(\sigma _t^2)={\displaystyle \frac{b_t}{a_t1}}=\widehat{\sigma }_S^2(r_t),`$
$`\text{Var}(\sigma _t^2)={\displaystyle \frac{b_t^2}{(a_t1)^2(a_t2)}}=2\widehat{\sigma }_S^4(r_t).`$
Solving the above equation, we obtain that
$$\widehat{a}_t=2.5\text{and}\widehat{b}_t=1.5\widehat{\sigma }_S^2(r_t).$$
Substituting this into (15), we obtain the following estimator
$$\widehat{\sigma }_{B,t}^2=\frac{1\lambda ^n}{1\lambda ^n+3(1\lambda )}\widehat{\sigma }_{ES,t}^2+\frac{3(1\lambda )}{1\lambda ^n+3(1\lambda )}\widehat{\sigma }_{S,t}^2.$$
(16)
Unfortunately, the weights in (16) are static, which does not depend on the time $`t`$. Hence, the Bayesian method does not produce a satisfactory answer to this problem.
## 5 Numerical Analysis
To facilitate the presentation, we use the simple abbreviation in Table 1 to denote five volatility estimation methods. Details of the first three methods can be found in Fan & Gu (2003). In particular, the first method is to estimate the volatility using the standard deviation of the yields in the past year and the RiskMetrics method is based on the exponential smoothing with $`\lambda =0.94`$. The semiparametric method of Fan & Gu (2003) is an extension of a local model used in the exponential smoothing, with the smoothing parameter determined by minimizing the prediction error. It includes the exponential smoothing with $`\lambda `$ selected by data as a specific example.
The following four measures are employed to assess the performance of different procedures for estimating the volatility. Other related measures can also be used. See Davé & Stahl (1997).
Measure 1. Exceedence ratio against confidence level.
This measure counts the number of the events for which the loss of an asset exceeds the loss predicted by the normal model at a given confidence $`\alpha `$. With estimated volatility, under the normal model, the one-period VaR is estimated by $`\mathrm{\Phi }^1(\alpha )\widehat{\sigma }_t`$, where $`\mathrm{\Phi }^1(\alpha )`$ is the $`\alpha `$ quantile of the standard normal distribution. For each estimated VaR, the Exceedence Ratio (ER) is computed as
$$\text{ER}(\widehat{\sigma }_t^2)=m^1\underset{i=T+1}{\overset{T+m}{}}I(Y_i<\mathrm{\Phi }^1(\alpha )\widehat{\sigma }_i),$$
(17)
for an out-sample of size $`m`$. This gives an indication on how effective the volatility estimator can be used for predicting the one-period VaR. Note that the Monte Carlo error for this measure has an approximate size $`\{\alpha (1\alpha )/m\}^{1/2}`$, even when the true $`\sigma _t`$ is used. For example, with $`\alpha =5\%`$ and $`m=1000`$, the Monte Carlo error is around $`0.68\%`$. Thus, unless the post-sample size $`m`$ is large enough, this measure has difficulty in differentiating the performance of various estimators due to the presence of large error margins. Note that the ER depends strongly on the assumption of normality. If the underlying return process is non-normal, the Student’s $`t(5)`$ say, the ER will grossly be overestimated even with the true volatility. In our simulation study, we will employ the true $`\alpha `$-quantile of the error distribution instead of $`\mathrm{\Phi }^1(\alpha )`$ in (17) to compute the ER. For real data analysis, we use the $`\alpha `$-quantile of the last $`250`$ residuals for the in-sample data.
Measure 2. Mean Absolute Deviation Error.
To motivate this measure, let us first consider the mean square errors:
$$\text{PE}(\widehat{\sigma }_t^2)=m^1\underset{i=T+1}{\overset{T+m}{}}(Y_i^2\widehat{\sigma }_i^2)^2.$$
The expected value can be decomposed as
$$E(\text{PE})=m^1\underset{i=T+1}{\overset{T+m}{}}E(\sigma _i^2\widehat{\sigma }_i^2)^2+m^1\underset{i=T+1}{\overset{T+m}{}}E(Y_i^2\sigma _i^2)^2.$$
(18)
Note that the first term reflects the effectiveness of the estimated volatility while the second term is the size of the stochastic error, independent of estimators. As in all statistical prediction problems, the second term is usually of an order of magnitude larger than the first term. Thus, a small improvement on PE could mean substantial improvement over the estimated volatility. However, due to the well-known fact that financial time series contain outliers, the mean-square error is not a robust measure. Therefore, we used the mean-absolute deviation error (MADE):
$$\text{MADE}(\widehat{\sigma }_t^2)=m^1\underset{i=T+1}{\overset{T+m}{}}Y_i^2\widehat{\sigma }_i^2.$$
Measure 3. Square-root Absolute Deviation Error.
An alternative variation to MADE is the square-Root Absolute Deviation Error (RADE), which is defined as
$$\text{RADE}(\widehat{\sigma }_t^2)=m^1\underset{i=T+1}{\overset{T+m}{}}\left|Y_i\sqrt{\frac{2}{\pi }}\widehat{\sigma }_i\right|.$$
The constant factor comes from the fact that $`E|\epsilon _t|=\sqrt{\frac{2}{\pi }}`$ for $`\epsilon _tN(0,1)`$. If the underlying error distribution deviates from normality, this measure is not robust.
Measure 4. Ideal Mean Absolute Deviation Error.
To assess the estimation of the volatility in simulations, one can also employ the ideal mean absolute deviation error (IMADE):
$$\text{IMADE}=m^1\underset{i=T+1}{\overset{T+m}{}}|\widehat{\sigma }_i^2\sigma _i^2|.$$
This measure calibrates the accuracy of the forecasted volatility in terms of the absolute difference between the true and the forecasted one. However, for real data analysis, this measure is not applicable.
### 5.1 Simulations
To assess the performance of the five estimation methods in Table 1, we compute the average and the standard deviation of each of the four measures over $`600`$ simulations. Generally speaking, the smaller the average (or the standard deviation), the better the estimation approach. We also compute the “score” of an estimator, which is the percentage of times among 600 simulations that the estimator outperforms the average of the 5 methods in terms of an effectiveness measure. To be more specific, for example, consider RiskMetrics using MADE as an effectiveness measure. Let $`m_i`$ be the MADE of the RiskMetrics estimator at the $`i`$-th simulation, and $`\overline{m}_i`$ the average of the MADEs for the five estimators at the $`i`$-th simulation. Then the “score” of the RiskMetrics approach in terms of the MADE is defined as
$$\frac{1}{600}\underset{i=1}{\overset{600}{}}I(m_i<\overline{m}_i).$$
Obviously, the estimators with higher scores are preferred. In addition, we define a “relative loss” of an estimator $`\widehat{\sigma }_t^2`$ relative to $`\widehat{\sigma }_{I,t}^2`$ in terms of MADEs as
$$\text{RLOSS}(\widehat{\sigma }_t^2\widehat{\sigma }_{I,t}^2)=\frac{\overline{\text{MADE}}(\widehat{\sigma }_t^2)\overline{\text{MADE}}(\widehat{\sigma }_{I,t}^2)}{\overline{\text{MADE}}(\widehat{\sigma }_{I,t}^2)},$$
where $`\overline{\text{MADE}}(\widehat{\sigma }_t^2)`$ is the average of MADE($`\widehat{\sigma }_t^2`$) among simulations.
Example 1. To simulate the interest rate data, we consider the Cox-Ingersoll-Ross (CIR) model:
$`dr_t=\kappa (\theta r_t)dt+\sigma r_t^{1/2}dW_t,tt_0,`$
where the spot rate, $`r_t`$, moves around a central location or long-run equilibrium level $`\theta =0.08571`$ at speed $`\kappa =0.21459`$. The $`\sigma `$ is set to be 0.07830. These values of parameters are cited from Chapman & Pearson (2000), which satisfy the condition $`2\kappa \theta \sigma ^2`$ so that the process $`r_t`$ is stationary and positive. The model has been studied by Chapman & Pearson (2000) and Fan & Zhang (2003).
There are two methods to generate samples from this model. The first one is the discrete-time order $`1.0`$ strong approximation scheme in Kloeden, et al. (1996); the second one is using the exact transition density detailed in Cox et al. (1985) and Fan & Zhang (2003). Here we use the first method to generate $`600`$ series of data each with length $`1200`$ of the weekly data from this model. For each simulation, we set the first $`900`$ observations as the “in-sample” data and the last $`300`$ observations as the “out-sample” data.
The results are summarized in Table 2, which shows that the performance of the integrated estimator uniformly dominates the other estimators because of its highest score, lowest IMADE, MADE, and RADE. The improvement in IMADE is over $`100`$ percent. This shows that our integrated volatility method better captures the volatility dynamics. The Bayesian method of combining the estimates from the time and state domains outperforms all other methods. The historical simulation method performed poorly due to mis-specification of the function of the volatility parameter. The results here show the advantage of aggregating the information of time domain and state domain. Note that all estimators have reasonable ER values at level $`0.05`$, especially the ER value of the integrated estimator is closest to $`0.05`$. To appreciate how much improvement for our integrated method over the other methods, we display the mean absolute difference between the forecasted and the true volatility in Figure 2. It is seen that the integrated method is much better than the others in terms of the difference.
Example 2. There is a large literature on the estimation of volatility. In addition to the famous parametric models such as ARCH and GARCH, stochastic volatility models have also received a lot of attention. For an overview, see, for example, Barndoff-Neilsen & Shephard (2001, 2002), Bollerslev & Zhou (2002) and references therein. We consider the following stochastic volatility model:
$`dr_t=\sigma _tdB_t,r_0=0`$
$`dV_t=\kappa (\theta V_t)dt+\alpha V_tdW_t,V_0=\eta ,V_t=\sigma _t^2,`$
where $`W_t`$ and $`B_t`$ are two independent standard Brownian motions.
There are two methods to generate samples from this model. One is the direct method, using the result of Genon-Catalot et al. (1999). Let $`a=1+2\kappa /\alpha ^2`$ and $`b=2\theta \kappa /\alpha ^2`$. The conditions (A1)-(A4) in the above paper are satisfied with the parameter values in the model being constants as $`\kappa =3`$, $`\theta =0.009`$ and $`\alpha ^2=4`$ and the initial random variable $`\eta `$ follows the Inverse Gamma distribution. The value of $`\theta `$ is set as the real variance of the daily return for Standard & Poor 500 data from January 4, 1988 to December 29, 2000. The value $`\alpha ^2`$ is to make the parameter $`a`$ of the stable distribution $`IG(a,b)`$ equal $`2.5`$, the prior parameter in (10). If $`\mathrm{\Delta }0`$ and $`n\mathrm{\Delta }\mathrm{}`$, then
$$Y_i\sqrt{\frac{b}{a}}T,\text{where}Tt(2a).$$
Another method is the discretization of the model. Conditionally on $`\text{g}=\sigma (V_t,t0)`$, the random variables $`Y_i`$ are independent and follow $`N(0,\overline{V}_i)`$ with
$$\overline{V}_i=\frac{1}{\mathrm{\Delta }}_{(i1)\mathrm{\Delta }}^{i\mathrm{\Delta }}V_s𝑑s.$$
To simulate the diffusion process $`V_t`$, one can use the following order 1.0 scheme with sampling interval $`\mathrm{\Delta }^{}=\mathrm{\Delta }/30`$,
$`V_{i+\mathrm{\Delta }^{}}=V_i+\kappa (\theta V_i)\mathrm{\Delta }^{}+\alpha V_i(\mathrm{\Delta }^{})^{1/2}\epsilon _i+{\displaystyle \frac{1}{2}}\alpha ^2V_i\mathrm{\Delta }^{}(\epsilon _i^21),`$
where $`\{\epsilon _i\}`$ are independent random series from the standard normal distribution.
We simulate $`600`$ series of $`1000`$ monthly data using the second method with step size $`\mathrm{\Delta }=1/12`$. For each simulated series, set the first three quarters observations as the in-sample data and the remaining observations as the out-sample data. The performance of each volatility estimation is described in Table 3. The conclusion similar to Example 1 can be drawn from this example.
Example 3. We now consider the geometric Brownian (GBM):
$$dr_t=\mu r_t+\sigma r_tdW_t,$$
where $`W_t`$ is a standard one-dimensional Brownian motion. This is a non-stationary process to which we check if our method continues to apply. Note that the celebrated Black-Scholes option price formula is derived on the Osborne’s assumption that the stock price follows the GBM model. By the It$`\widehat{o}`$ formula, we have
$$\mathrm{log}r_t\mathrm{log}r_0=(\mu \sigma ^2/2)t+\sigma ^2W_t.$$
We set $`\mu =0.03`$ and $`\sigma =0.26`$ in our simulations. With the Brownian motion simulated from independent Gaussian increments, one can generate the samples for the GBM. Here we use the latter with $`\mathrm{\Delta }=1/52`$ in $`600`$ simulations. For each simulation, we generate $`1000`$ observations and use the first two thirds of observations as in-sample data and the remaining observations as out-sample data.
Table 4 summarizes the results. The historical simulation approach has the smallest MADE, but suffers from poor forecast in terms of IMADE. This is surprising. Why is it so different between IMADE and MADE? This phenomenon may be produced by the non-stationarity of the process. For the integrated method, even though the true volatility structure is well captured because of the lowest IMADE, extreme values of observations make the MADE quite large. To more accurately calibrate the performance of the volatility estimation, we use the $`95\%`$ up-trimmed mean instead of the mean to summarize the values of the measures. Table 5 reports the trimmed means and the relative losses for different measures. The similar conclusions to those in Example 1 can be drawn from the table. This shows that our integrated method continues to perform better than other for this non-stationary case. The Bayesian estimator performs comparably with the dynamically integrated method and outperforms all others.
### 5.2 Empirical Study
In this section, we will apply the integrated volatility estimation methods and others to the analysis of real financial data.
#### 5.2.1 Treasury Bond
We consider here the weekly returns of three treasury bonds with terms 1, 5 and 10 years, respectively.
We set the observations from January 4, 1974 to December 30, 1994 as in-sample data, and those from January 6, 1995 up to August 8, 2003 as out-sample data. The total sample size is $`1545`$ and the in-sample size is $`1096`$. The results are reported in Table 6.
From Table 6, the integrated estimator is of the smallest MADE and almost the smallest RADE, which reflects that the integrated estimation method of the volatility is the best among the five methods. Relative losses in MADE of the other estimators with respect to the integrated estimator can easily be computed as ranging from $`8.47\%`$ (NonBay) to $`42.6\%`$ (Hist) for the bond with one year term. For the bonds with 5 or 10 years term, the five estimators have close MADEs and RADEs, where the historical simulation method is better than the RiskMetrics in terms of MADE and RADE, and the integrated estimation approach has the smallest MADEs. This demonstrates the advantage of using state domain information which can help the time-domain prediction of the changes in bond interest dynamics.
#### 5.2.2 Exchange Rate
We analyse the daily exchange rate of several foreign currencies with US dollar. The data are from January 3, 1994 to August 1, 2003. The in-sample data consists of the observations before January 1, 2001, and the out-sample data consists of the remaining observations. The results are reported in Table 7. It is seen that the integrated estimator has the smallest MADEs for the exchange rates, which again supports our integrated volatility estimation.
## 6 Conclusions
We have proposed a Bayesian method and a dynamically integrated method to aggregate the information from the time-domain and the state domain. The performance comparisons are studied both empirically and theoretically. We have shown that the proposed integrated method is effectively aggregating the information from both the time and the state domains, and has advantages over some previous methods. It is powerful in forecasting volatilities for the yields of bonds and for exchange rates. Our study has also revealed that proper use of information from both the time domain and the state domain makes volatility forecasting more accurately. Our method exploits the continuity in the time-domain and stationarity in the state-domain. It can be applies to situations where these two conditions hold approximately.
## 7 Appendix
We collect technical conditions for the proof of our results.
$`\sigma ^2(x)`$ is Lipschitz continuous.
There exists a constant $`L>0`$ such that $`E|\mu (r_s)|^{2(p+\delta )}L`$ and $`E|\sigma (r_s)|^{2(p+\delta )}L`$ for any $`s[t\eta ,t]`$, where $`\eta `$ is some positive constant, $`p`$ is an integer not less than $`4`$ and $`\delta >0`$.
The discrete observations $`\{r_{t_i}\}_{i=0}^N`$ satisfy the stationarity conditions of Banon (1978). Furthermore, the $`G_2`$ condition of Rosenblatt (1970) holds for the transition operator.
The conditional density $`p_{\mathrm{}}(y|x)`$ of $`r_{t_{i+\mathrm{}}}`$ given $`r_{t_i}`$ is continuous in the arguments $`(y,x)`$ and is bounded by a constant independent of $`\mathrm{}`$.
The kernel $`W`$ is a bounded, symmetric probability density function with compact support, $`[1,1]`$ say.
$`(Nn)h\mathrm{}`$, $`(Nn)h^50`$, $`(Nn)h\mathrm{\Delta }0`$.
Throughout the proof, we denote by $`M`$ a generic positive constant, and use $`\mu _s`$ and $`\sigma _s`$ to represent $`\mu (r_s)`$ and $`\sigma (r_s)`$, respectively.
Proof of Proposition 1. It suffices to show that the process $`\{r_s\}`$ is Hölder-continuous with order $`q=(p1)/(2p)`$ and coefficient $`K_1`$, where $`E[K_1^{2(p+\delta )}]<\mathrm{}`$, because this together with assumption $`(A1)`$ gives the result of the lemma. By Jensen’s inequality and martingale moment inequalities (Karatzas & Shreve 1991, Section 3.3.D), we have
$`E|r_ur_s|^{2(p+\delta )}`$ $`M\left(E\left|{\displaystyle _s^u}\mu _v𝑑v\right|^{2(p+\delta )}+E\left|{\displaystyle _s^u}\sigma _v𝑑W_v\right|^{2(p+\delta )}\right)`$
$`M(us)^{2(p+\delta )1}{\displaystyle _s^u}E|\mu _v)|^{2(p+\delta )}dv+M(us)^{p+\delta 1}{\displaystyle }_s^uE|\sigma _v|^{2(p+\delta )}dv`$
$`M(us)^{p+\delta }.`$
Then by the Kolmogorov continuity theorem (Revuz & Yor 1991, Theorem 2.1), $`\{r_s\}`$ is Hölder-continuous.
Proof of Theorem 1. Let $`Z_{i,s}=(r_sr_{t_i})^2`$. Applying Itô formula to $`Z_{i,s}`$, we obtain
$`dZ_{i,s}=`$ $`2\left({\displaystyle _{t_i}^s}\mu _u𝑑u+{\displaystyle _{t_i}^s}\sigma _u𝑑W_u\right)\left(\mu _sds+\sigma _sdW_s\right)+\sigma _s^2ds`$
$`=`$ $`2\left[\left({\displaystyle _{t_i}^s}\mu _u𝑑u+{\displaystyle _{t_i}^s}\sigma _u𝑑W_u\right)\mu _sds+\sigma _s\left({\displaystyle _{t_i}^s}\mu _u𝑑u\right)dW_s\right]`$
$`+2\left({\displaystyle _{t_i}^s}\sigma _u𝑑W_u\right)\sigma _sdW_s+\sigma _s^2ds.`$
Then $`Y_i^2`$ can be decomposed as
$$Y_i^2=2a_i+2b_i+\overline{\sigma }_i^2,$$
where
$$a_i=\mathrm{\Delta }^1\left[_{t_i}^{t_{i+1}}\mu _s𝑑s_{t_i}^s\mu _u𝑑u+_{t_i}^{t_{i+1}}\mu _s𝑑s_{t_i}^s\sigma _u𝑑W_u+_{t_i}^{t_{i+1}}\sigma _s𝑑W_s_{t_i}^s\mu _u𝑑u\right],$$
$$b_i=\mathrm{\Delta }^1_{t_i}^{t_{i+1}}_{t_i}^s\sigma _u𝑑W_u\sigma _s𝑑W_s,$$
and
$$\overline{\sigma }_i^2=\mathrm{\Delta }^1_{t_i}^{t_{i+1}}\sigma _s^2𝑑s.$$
Therefore, $`\widehat{\sigma }_{ES,t}^2`$ can be written as
$`\widehat{\sigma }_{ES,t}^2`$ $`=`$ $`2{\displaystyle \frac{1\lambda }{1\lambda ^n}}{\displaystyle \underset{i=tn}{\overset{t1}{}}}\lambda ^{ti1}a_i+2{\displaystyle \frac{1\lambda }{1\lambda ^n}}{\displaystyle \underset{i=tn}{\overset{t1}{}}}\lambda ^{ti1}b_i+{\displaystyle \frac{1\lambda }{1\lambda ^n}}{\displaystyle \underset{i=tn}{\overset{t1}{}}}\lambda ^{ti1}\overline{\sigma }_i^2`$
$``$ $`A_{n,\mathrm{\Delta }}+B_{n,\mathrm{\Delta }}+C_{n,\mathrm{\Delta }}.`$
By Proposition 1, as $`n\mathrm{\Delta }0`$,
$$|C_{n,\mathrm{\Delta }}\sigma _t^2|K(n\mathrm{\Delta })^q,$$
where $`q=(p1)/(2p)`$. This combined with Lemmas 1-2 below completes the proof of the theorem.
###### Lemma 1
If condition (A2) is satisfied, then $`E[A_{n,\mathrm{\Delta }}^2]=O(\mathrm{\Delta }).`$
Proof . Simple algrbea gives the result. In fact,
$`E(a_i^2)`$ $``$ $`3E\left[\mathrm{\Delta }^1{\displaystyle _{t_i}^{t_{i+1}}}\mu _s𝑑s{\displaystyle _{t_i}^s}\mu _u𝑑u\right]^2+3E\left[\mathrm{\Delta }^1{\displaystyle _{t_i}^{t_{i+1}}}\mu _s𝑑s{\displaystyle _{t_i}^s}\sigma _u𝑑W_u\right]^2`$
$`+3E\left[\mathrm{\Delta }^1{\displaystyle _{t_i}^{t_{i+1}}}\sigma _s𝑑W_s{\displaystyle _{t_i}^s}\mu _u𝑑u\right]^2`$
$``$ $`I_1(\mathrm{\Delta })+I_2(\mathrm{\Delta })+I_3(\mathrm{\Delta }).`$
Applying Jensen’s inequality, we obtain that
$`I_1(\mathrm{\Delta })`$ $`=`$ $`O(\mathrm{\Delta }^1)E\left[{\displaystyle _{t_i}^{t_{i+1}}}{\displaystyle _{t_i}^s}\mu _s^2\mu _u^2𝑑u𝑑s\right]`$
$`=`$ $`O(\mathrm{\Delta }^1){\displaystyle _{t_i}^{t_{i+1}}}{\displaystyle _{t_i}^s}E(\mu _u^4+\mu _s^4)𝑑u𝑑s=O(\mathrm{\Delta }).`$
By Jensen’s inequality, Hölder’s inequality and martingale moments inequalities, we have
$`I_2(\mathrm{\Delta })`$ $`=`$ $`O(\mathrm{\Delta }^1){\displaystyle _{t_i}^{t_{i+1}}}E\left(\mu _s{\displaystyle _{t_i}^s}\sigma _u^2𝑑W_u\right)^2𝑑s`$
$`=`$ $`O(\mathrm{\Delta }^1){\displaystyle _{t_i}^{t_{i+1}}}\left\{E\left[\mu _s\right]^4E\left[{\displaystyle _{t_i}^{t_{i+1}}}\sigma _u𝑑W_u\right]^4\right\}^{1/2}𝑑s=O(\mathrm{\Delta }).`$
Similarly, $`I_3(\mathrm{\Delta })=O(\mathrm{\Delta }).`$ Therefore, $`E(a_i^2)=O(\mathrm{\Delta })`$. Then by the Cauchy-Schwartz inequality and noting that $`n(1\lambda )=O(1)`$, we obtain that
$$E[A_{n,\mathrm{\Delta }}^2]n\left(\frac{1\lambda }{1\lambda ^n}\right)^2\underset{i=1}{\overset{n}{}}\lambda ^{2(ni)}E(a_i^2)=O(\mathrm{\Delta }).$$
###### Lemma 2
Under condition (A2), if $`n\mathrm{}`$ and $`n\mathrm{\Delta }0`$, then
$$s_{1,t}^1\sqrt{n}B_{n,\mathrm{\Delta }}\stackrel{𝒟}{}𝒩(0,1).$$
(A1)
Proof. Note that
$$b_j=\sigma _t^2\mathrm{\Delta }^1_{t_j}^{t_{j+1}}(W_sW_{t_j})𝑑W_s+ϵ_j,$$
where
$$ϵ_j=\mathrm{\Delta }^1_{t_j}^{t_{j+1}}(\sigma _s\sigma _t)\left[_{t_j}^s\sigma _u𝑑W_u\right]𝑑W_s+\mathrm{\Delta }^1\sigma _t_{t_j}^{t_{j+1}}\left[_{t_j}^s(\sigma _u\sigma _t)𝑑W_u\right]𝑑W_s.$$
By the central limit theorem for martingale (see Hall & Heyde 1980, Corollary 3.1), it suffices to show that
$$V_n^2E[s_{1,t}^2nB_{n,\mathrm{\Delta }}^2]1,$$
(A2)
and the following Lyapunov condition holds:
$$\underset{i=tn}{\overset{t1}{}}E\left(\sqrt{n}\frac{1\lambda }{1\lambda ^n}\lambda ^{ti1}b_i\right)^40.$$
(A3)
Note that
$`{\displaystyle \frac{\mathrm{\Delta }^2}{2}}E(ϵ_j^2)`$ $``$ $`E\left\{{\displaystyle _{t_j}^{t_{j+1}}}(\sigma _s\sigma _t)\left[{\displaystyle _{t_j}^s}\sigma _u𝑑W_u\right]𝑑W_s\right\}^2`$ (A4)
$`+\sigma _t^2E\left\{{\displaystyle _{t_j}^{t_{j+1}}}\left[{\displaystyle _{t_j}^s}(\sigma _u\sigma _{t_t})𝑑W_u\right]𝑑W_s\right\}^2`$
$``$ $`L_{n1}+L_{n2}.`$
By Jensen’s inequality, Hölder’s inequality and moments inequalities for martingale, we have
$`L_{n1}`$ $``$ $`{\displaystyle _{t_j}^{t_{j+1}}}E\left\{(\sigma _s\sigma _t)^2\left[{\displaystyle _{t_j}^s}\sigma _u𝑑W_u\right]^2\right\}𝑑s`$ (A5)
$``$ $`{\displaystyle _{t_j}^{t_{j+1}}}\left\{E(\sigma _s\sigma _t)^4E\left[{\displaystyle _{t_j}^s}\sigma _u𝑑W_u\right]^4\right\}^{1/2}𝑑s`$
$``$ $`{\displaystyle _{t_j}^{t_{j+1}}}\left\{E[K(n\mathrm{\Delta })^q]^4\mathrm{\hspace{0.17em}36}\mathrm{\Delta }{\displaystyle _{t_j}^s}E(\sigma _u^4)𝑑u\right\}^{1/2}𝑑s`$
$``$ $`M(n\mathrm{\Delta })^{2q}\mathrm{\Delta }^2.`$
Similarly,
$$L_{n2}M(n\mathrm{\Delta })^{2q}\mathrm{\Delta }^2.$$
(A6)
By (A4), (A5) and (A6),
$$E(ϵ_j^2)M(n\mathrm{\Delta })^{2q}.$$
(A7)
Therefore,
$$E[\sigma _t^4b_j^2]=\frac{1}{2}+O((n\mathrm{\Delta })^q).$$
By the theory of stochastic calculus, simple algebra gives that $`E(b_j)=0`$ and $`E(b_ib_j)=0`$ for $`ij`$. It follows that
$$V_n^2=E(s_{1,t}^2nB_{n,\mathrm{\Delta }}^2)=\underset{i=tn}{\overset{t1}{}}E\left(2s_{1,t}\sqrt{n}\frac{1\lambda }{1\lambda ^n}\lambda ^{ti1}b_i\right)^21.$$
That is, (A2) holds. For (A3), it suffices to prove that $`E(b_j^4)`$ is bounded, which holds by applying the moment inequalities for martingales to $`b_j^4`$.
Proof of Theorem 2. The proof is completed by using the same lines in Fan & Zhang (2003).
Proof of Theorem 3. By Fan & Yao (1998), the volatility estimator $`\widehat{\sigma }_{S,t_N}^2`$ behaves as if the instantaneous return function $`f`$ is known, hence without loss of generality we assume that $`f(x)=0`$ and hence $`\widehat{R}_i=Y_i^2.`$ Let $`𝐘=(Y_0^2,\mathrm{},Y_{Nn1}^2)^T`$, $`𝐖=\text{diag}\{W_h(r_{t_0}r_{t_N}),\mathrm{},W_h(r_{t_{Nn1}}r_{t_N})\},`$ and
$$𝐗=\left(\begin{array}{cc}1& r_{t_0}r_{t_N}\\ \mathrm{}& \mathrm{}\\ 1& r_{t_{Nn1}}r_{t_N}\end{array}\right).$$
Denote by $`m_i=E[Y_i^2|r_{t_i}]`$, $`𝐦=(m_0,\mathrm{},m_{Nn1})^T`$ and $`𝐞_\mathrm{𝟏}=(1,0)^T`$. Define $`𝐒_𝐍=𝐗^𝐓\mathrm{𝐖𝐗}`$ and $`𝐓_𝐍=𝐗^𝐓\mathrm{𝐖𝐘}`$. Then it can be written that (see Fan & Yao, 2003)
$$\widehat{\sigma }_{S,t_N}^2=𝐞_\mathrm{𝟏}^𝐓𝐒_𝐍^\mathrm{𝟏}𝐓_𝐍.$$
Hence
$`\widehat{\sigma }_{S,t_N}^2\sigma _{t_N}^2`$ $`=`$ $`𝐞_\mathrm{𝟏}^𝐓𝐒_𝐍^\mathrm{𝟏}𝐗^𝐓𝐖\{𝐦𝐗𝜷_N\}+𝐞_\mathrm{𝟏}^𝐓𝐒_𝐍^\mathrm{𝟏}𝐗^𝐓𝐖(𝐘𝐦)`$ (A8)
$``$ $`𝐞_\mathrm{𝟏}^𝐓𝐛+𝐞_\mathrm{𝟏}^𝐓𝐭,`$
where $`𝜷_N=(m(r_{t_N}),m^{}(r_{t_N}))^T`$ with $`m(r_{t_N})=E[Y_1^2|r_{t_1}=r_{t_N}]`$. By Fan & Zhang (2003), the bias vector $`𝐛`$ converges in probability to a vector $`\overline{𝐛}`$ with $`\overline{𝐛}=O(h^2)=o(1/\sqrt{(Nn)h})`$. In the following, we will show that the centralized vector $`𝐭`$ is asymptotically normal.
In fact, put $`𝐮=(Nn)^1𝐇^\mathrm{𝟏}𝐗^𝐓𝐖(𝐘𝐦)`$ where $`𝐇=\text{diag}\{1,h\}`$, then by Fan & Zhang (2003) the vector $`𝐭`$ can be written as
$$𝐭=p^1(r_{t_N})𝐇^\mathrm{𝟏}𝐒^\mathrm{𝟏}𝐮(1+o_p(1)),$$
(A9)
where $`𝐒=(\mu _{i+j2})_{i,j=1,2}`$ with $`\mu _j=u^jW(u)𝑑u`$. For any constant vector $`𝐜`$, define
$$Q_N=𝐜^𝐓𝐮=\frac{1}{Nn}\underset{i=0}{\overset{Nn1}{}}\{Y_i^2m_i\}C_h(r_{t_i}r_{t_N}),$$
where $`C_h()=1/hC(/h)`$ with $`C(x)=c_0W(x)+c_1xW(x)`$. Applying the “big-block” and “small-block” arguments in Fan & Yao (2003, Theorem 6.3), we obtain
$$\theta ^1(r_{t_N})\sqrt{(Nn)h}Q_N\stackrel{\text{D}}{}N(0,1),$$
(A10)
where $`\theta ^2(r_{t_N})=2p(r_{t_N})\sigma ^4(r_{t_N})_{\mathrm{}}^+\mathrm{}C^2(u)𝑑u`$. In the following, we will decompose $`Q_N`$ into two parts, $`Q_N^{}`$ and $`Q_N^{\prime \prime }`$, which satisfy that
* $`(Nn)hE[\theta ^1(r_{t_N})Q_N^{}]^2\frac{h}{Nn}\left(h^1a_N(1+o(1))+(Nn)o(h^1)\right)0.`$
* $`Q_N^{\prime \prime }`$ is identically distributed as $`Q_N`$ and is asymptotically independent of $`\widehat{\sigma }_{ES,t_N}^2`$.
Define
$$Q_N^{}=\frac{1}{Nn}\underset{i=0}{\overset{a_N}{}}\{Y_i^2E[Y_i^2|r_{t_i}]\}C_h(r_{t_i}r_{t_N}),$$
(A11)
and
$$Q_N^{\prime \prime }=Q_NQ_N^{},$$
where $`a_N`$ is a positive integer satisfying $`a_N=o(Nn)`$ and $`a_N\mathrm{\Delta }\mathrm{}`$. Let $`\vartheta _{N,\mathrm{}}=(Y_i^2m_i)C_h(r_{t_i}r_{t_N})`$, then by Fan & Zhang (2003)
$$\text{Var}[\theta ^1(r_{t_N})\vartheta _{N,1}]=h^1(1+o(1))\text{ and }\underset{\mathrm{}=1}{\overset{Nn2}{}}|\text{Cov}(\vartheta _{N,1},\vartheta _{N,\mathrm{}+1})|=o(h^1),$$
(A12)
which yields the result in (i). This combined with (A10), (i) and (A11) leads to
$$\theta ^1(r_{t_N})\sqrt{(Nn)h}Q_N^{\prime \prime }\stackrel{\text{D}}{}N(0,1).$$
(A13)
Note that the stationarity conditions of Banon (1978) and the $`G_2`$ condition of Rosenblatt (1970) on the transition operator imply that the $`\rho `$-mixing coefficient $`\rho (\mathrm{})`$ of $`\{r_{t_i}\}`$ decays exponentially, and the strong-mixing coefficient $`\alpha (\mathrm{})\rho (\mathrm{})`$, it follows that
$$|E\mathrm{exp}\{i\xi (Q_N^{\prime \prime }+\widehat{\sigma }_{ES,t_N}^2)\}E\mathrm{exp}\{i\xi (Q_N^{\prime \prime }\}E\mathrm{exp}\{i\xi \widehat{\sigma }_{ES,t_N}^2)\}|32\alpha (s_N)0,$$
(A14)
for any $`\xi `$. Using the theorem of Volkonskii & Rozanov (1959), one gets the asymptotic independence of $`\widehat{\sigma }_{ES,t_N}^2`$ and $`Q_N^{\prime \prime }`$.
By (i), $`\sqrt{(Nn)h}Q_N^{}`$ is asymptotically negligible. This together with Theorem 1 lead to
$$d_1\theta ^1(r_{t_N})\sqrt{(Nn)h}Q_N+d_2V_2^{1/2}\sqrt{n}[\widehat{\sigma }_{ES,t_N}^2\sigma ^2(r_{t_N})]\stackrel{𝒟}{}𝒩(0,d_1^2+d_2^2),$$
for any $`d_1,d_2`$, where $`V_2=\frac{e^c+1}{e^c1}\sigma ^4(r_{t_N})`$. Since $`Q_N`$ is a linear transform of $`𝐮`$,
$`𝐕^{1/2}\left[\begin{array}{c}\sqrt{(Nn)h}𝐮\\ \sqrt{n}[\widehat{\sigma }_{ES,t_N}^2\sigma ^2(r_{t_N})]\end{array}\right]\stackrel{𝒟}{}𝒩(0,I_3),`$
where $`𝐕=\text{blockdiag}\{V_1,V_2\}`$ with $`V_1=2\sigma ^4(r_{t_N})p(r_{t_N})𝐒^{}`$, where $`𝐒^{}=(\nu _{i+j2})_{i,j=1,2}`$ with $`\nu _j=u^jW^2(u)𝑑u`$. This combined with (A9) gives the joint asymptotic normality of $`𝐭`$ and $`\widehat{\sigma }_{ES,t_N}^2`$. Note that $`𝐛=o_p(1/\sqrt{(Nn)h})`$, it follows that
$`\mathrm{\Sigma }^{1/2}\left(\begin{array}{c}\sqrt{(Nn)h}[\widehat{\sigma }_{S,t_N}^2\sigma ^2(r_{t_N})]\\ \sqrt{n}[\widehat{\sigma }_{ES,t_N}^2\sigma ^2(r_{t_N})]\end{array}\right)\stackrel{𝒟}{}𝒩(0,I_2),`$
where $`\mathrm{\Sigma }=\text{diag}\{2\sigma ^4(r_{t_N})\nu _0/p(r_{t_N}),V_2\}`$. Note that $`\widehat{\sigma }_{S,t_N}^2`$ and $`\widehat{\sigma }_{ES,t_N}^2`$ are asymptotically independent, it follows that the asymptotical normality of $`\widehat{\sigma }_{I,t_N}^2`$ holds.
Acknowledgements
The work was partially supported by a grant from the Research Grants Council of the Hong Kong SAR (Project No. CUHK 400903/03P), the NSF grant DMS-0355179 and the Chinese NSF grants 10471006 and 10001004. The authors thank Dr. Juan Gu for various assistances.
References
* Banon, G. (1978). Nonparametric identification for diffusion processes. SIAM J. Control Optim 16, 380-395.
* Barndoff-Neilsen, O.E. & Shephard, N. (2001). Non-Gaussian Ornstein-Uhlenbeck-based models and some of their uses in financial economics (with discussion. J. R. Statist. Soc. B 63, 167-241.
* Barndoff-Neilsen, O.E. & Shephard, N. (2002). Econometric analysis of realized volatility and its use in estimating stochastic volatility models. J. R. Statist. Soc. B 64, 253-280.
* Bollerslev, T. & Zhou, H. (2002). Estimating stochastic volatility diffusion using conditional moments of integrated volatility. Jour. Econometrics 109, 33-65.
* Chan, K.C., Karolyi, A.G., Longstaff, F.A. & Sanders, A.B. (1992). An empirical comparison of alternative models of the short-term interest rate. Journal of Finance 47, 1209-1227.
* Chapman, D.A. & Pearson, N.D. (2000). Is the short rate drift actually nonlinear? Journal of Finance 55, 355–388.
* Cox, J.C., Ingersoll, J.E. & Ross, S. A. (1985). A theory of the term structure of interest rates. Econometrica 53, 385-467.
* Davé, R. D. & Stahl, G. (1997). On the accuracy of VaR estimates based on the Variance-Covariance approach. Working paper, Olshen & Associates.
* Duffie, D. & Pan, J. (1997). An overview of Value at Risk. The Journal of Derivatives, 7–49.
* Fan, J. & Gu, J. (2003). Semiparametric estimation of value-at-risk. Econometrics Journal 6, 261-290.
* Fan, J. & Yao, Q. (1998). Efficient estimation of conditional variance functions in stochastic regression. Biometrika 85, 645–660.
* Fan, J. & Yao, Q. (2003). Nonlinear Time Series: Nonparametric and Parametric Methods, Springer-Verlag, New York.
* Fan, J. & Zhang, C.M. (2003). A Reexamination of Diffusion Estimators with Applications to Financial Model Validation. J. Am. Statist. Assoc. 98, 118–134.
* Genon-Catalot, Jeanthheau, T. & Laredo, C. (1999). Parameter estimation for discretely observed stochastic volatility models, Bernoulli 5, 855–872.
* Gijbels, I., Pope, A., & Wand, M.P. (1999). Understanding exponential smoothing via kernel regression. J. R. Statist. Soc. B 61, 39–50.
* Hall, P. & Heyde, C. (1980). Martingale limit theorem and its applications. Academic Press.
* Kloeden, D.E., Platen, E., Schurz, H. & Sørensen, M. (1996). On effects of discretization on estimators of drift parameters for diffusion processes. Journal of Applied Probability 33, 1061–1076.
* Karatzas, I. & Shreve, S. (1991). Brownian motion and stochastic calculus (2nd edition). Springer-Verlag, New York.
* Morgan, J.P. (1996) RiskMetrics Technical Document, Fourth edition, New York.
* Stanton, R. (1997). A nonparametric models of term structure dynamics and the market price of interest rate risk. Journal of Finance LII, 1973–2002.
* Revuz, D. & Yor, M. (1991). Continuous Martingales and Brownian Motion. Springer-Verlag.
* Rosenblatt, M. (1970). Density estimates and Markov sequences. In Nonparametric Techniques in Statistical Inference (ML Puri, ed.) 199-213. Cambridge Univ. Press.
* Ruppert, D., Wand, M.P., Holst, U. & Hössjer, O. (1997). Local polynomial variance function estimation. Technometrics 39, 262-273.
* Spokoiny, V. (2000). Drift estimation for nonparametric diffusion. Ann. Statist. 28, 815–836.
* Tong, H. (1990). Non-Linear Time Series: A Dynamical System Approach. Oxford University Press, Oxford.
* Tong, H. (1995). A personal overview of non-linear time series analysis from a chaos perspective (with discussion). Scandinavian Journal of Statistics 22, 399-445.
* Zhang, C.M. (2003). Calibrating the degrees of freedom for automatic data-smoothing and effective curve-checking. J. Am. Statist. Assoc. 98, 609-628 . |
warning/0506/cond-mat0506537.html | ar5iv | text | # Density functional theory for a macroion suspension
## I Introduction
The phase behaviour of charged colloidal suspensions at low ionic strength has attracted much experimental and theoretical interest. Observations of void structures and other phenomena phe ; Gröhn and Antonietti (2000) motivated a number of theoretical studies which attributed the anomalous behaviour to phase separation between colloid-rich and colloid-poor phases the ; van Roij et al. (1999); Warren (2000); Chan (2001). Several reviews are available rev ; Belloni (2000). The original theoretical explanations have come under strong attack for using a Debye-Hückel linearisation approximation which is, at best, at the margin of its validity. Various attempts to patch this up have left the situation unclear. Cell model calculations using Poisson-Boltzmann theory indicate the original predictions are an artefact of the linearisation approximation GvR ; Tamashiro and Schiessel (2003). The Debye-Hückel approximation can be improved by taking into account counterion condensation, in which case the phase transition may or may not be recovered depending on the approximation scheme used dhf . Other approaches such as extended Debye-Hückel theory Chan et al. (2001), symmetrised Poisson-Boltzmann theory Bhuiyan and Outhwaite (2002, 2004), ‘boot-strap’ Poisson-Boltzmann theory Petris et al. (2003), and a systematic expansion into two- and three-body interactions thr ; Hynninen et al. (2004), all indicate that a phase transition can occur, as do several integral equation studies int ; González-Tovar (1999). The experimental situation is also uncertain since a plausible alternative explanation has been suggested Belloni (2000), in which the voids correspond to regions occupied by dilute, highly extended (and therefore effectively invisble) polyelectrolyte chains which have been shed by the latex colloids.
Simulation methods struggle to approach these problems because of the disparity in size between the macroions and the small ions, and the need to handle the electrostatic interactions. Nevertheless, convincing evidence has been found for liquid-liquid phase separation in a macroion system at lower dimensionless temperatures sim ; Reščič and Linse (2001). Experimentally this corresponds to a solvent with a lower dielectric constant than water (but one in which the ions still disperse). Charged colloids in such solvents exhibit many interesting phenomena non .
Thus, whilst the weight of evidence perhaps suggests that aqueous charge-stabilised colloidal suspensions may not show genuine liquid-liquid phase coexistence, it is absolutely clear that there will be phase coexistence between condensed and dilute colloidal phases at small enough dimensionless temperatures. In this sense, the problem resembles the much-studied restricted primitive model (RPM), whose phase behaviour is now well established Fisher (1994); Stell (1995); Luijten et al. (2002); Kim et al. (2003).
In situations where genuine phase coexistence obtains, one can go on to ask questions about the surface tension and electrical structure of the interface between the coexisting phases. Answers to these questions may prompt new avenues for experimental investigation of real systems. Previously, Knott and Ford compute the surface tension using square-gradient theory, but discard the possible electrical structure at the interface Knott and Ford (2002). The present work approaches this problem within the context of a density functional theory, motivated by the earlier study in Ref. Warren (2000) (see also Appendix A). It places the phenomenological remarks made in this earlier work on a sounder footing. The analysis in Ref. Warren (2000) suggests that the macroion self energy is the dominant contribution to the excess free energy, similar to an early insight by Langmuir Langmuir (1938). In the present work therefore, the rather gross simplification has been adopted in which the macroion self energy is the *only* contribution to the excess free energy. Moreover this self energy is computed in a simple closed form using Debye-Hückel theory, and is thus also based on the much-criticised linearisation approximation. Nevertheless I argue that it is instructive to proceed, because of the rich phenomenology that is revealed.
The model predicts phase separation at low dimensionless temperatures and low ionic strengths, and in quantitative terms stands reasonable comparison with some of the other approaches. The physics of the phase separation lies in the dependence of the macroion self energy on the local ionic strength: macroions drift towards regions of high ionic strength, which by charge neutrality are regions where other macroions have also congregated. Within the linearisation approximation, the effect grows without bound as the macroion charge is increased, and thus the mechanism can drive phase separation at sufficiently large macroion charges. In reality, non-linear effects (counterion condensation) limit the effective macroion charge sat , and therefore this mechanism is probably insufficient in itself to drive phase separation in real systems. Undoubtably though it is still a contributing factor, operating in conjunction with other effects such as correlated fluctuations in the counterion clouds around macroions and the sharing of counterions between macroions Schmitz (1997); Gröhn and Antonietti (2000); sim .
The model is constructed in the form of a density functional theory. Thus, as well as making predictions for phase separation, it can be used to solve for the density profiles and the surface tension between coexisting phases. The results obtained here are in accord with typical expectations for soft condensed matter systems sur , and were summarised in an earlier publication Warren (2003). In addition, I also discuss the predictions that the theory makes for the structure factors. These are found to obey the Stillinger-Lovett moment conditions Stillinger and Lovett (1968); Martin (1988), although it turns out this is not a stringent test of the theory. Intriguingly, I find that the structure factors may diverge at a non-zero wavevector as one approaches the critical points. This suggests the possibility that the critical points in these systems may be replaced by charge-density-wave phases War . This phenomenological possibility in charged systems was first suggested by Nabutovskii and coworkers Nabutovskii et al. (1980); Hoye and Stell (1990); Fisher (1994).
## II Specification of the model
The underlying model of the macroion suspension used here is a primitive model commonly deployed for this kind of problem. The ‘primitive’ aspect is that the solvent is treated as a structureless dielectric continuum in which the macroions and small salt ions are embedded. The macroions are treated as spheres of (positive) charge $`Z`$, diameter $`\sigma `$, and number density $`\rho _m`$ (volume fraction $`\varphi =\pi \sigma ^3\rho _m/6`$). The salt ions are univalent counterions and coions at number densities $`\rho _{}`$ and $`\rho _+`$ respectively. I suppose there is only one species of counterion. The size of the salt ions is assumed to be small enough to be irrelevant. The dielectric continuum is characterised by a Bjerrum length $`l_\mathrm{B}`$ so that the electrostatic interaction energy between a pair of univalent charges separated by a distance $`r`$ is $`l_\mathrm{B}/r`$, in units of $`k_\mathrm{B}T`$ where $`k_\mathrm{B}`$ is Boltzmann’s constant and $`T`$ is the temperature. For water at room temperature, $`l_\mathrm{B}0.72\mathrm{nm}`$. The model is completely parametrised by the dimensionless ratio $`\sigma /l_\mathrm{B}`$ and the charge $`Z`$. It is often convenient to pretend that the dielectric permittivity of the background is independent of temperature, in which case $`l_\mathrm{B}1/T`$. This means that $`\sigma /l_\mathrm{B}`$ can be regarded as a dimensionless temperature.
The density functional theory is specified by giving the free energy $`F`$ as a functional of the spatially varying number densities $`\rho _m(𝐫)`$ and $`\rho _\pm (𝐫)`$ Eva . The functional is decomposed into ideal, mean-field, and correlation contributions:
$$\begin{array}{c}\frac{F}{k_\mathrm{B}T}=d^3𝐫\underset{i=m,\pm }{}\rho _i(𝐫)\mathrm{ln}\frac{\rho _i(𝐫)}{e\rho _i^{}}\hfill \\ +\frac{l_\mathrm{B}}{2}d^3𝐫d^3𝐫^{}\frac{\rho _z(𝐫)\rho _z(𝐫^{})}{|𝐫𝐫^{}|},\hfill \\ +\frac{1}{k_\mathrm{B}T}d^3𝐫\rho _m(𝐫)f_m(𝐫).\hfill \end{array}$$
(1)
The first term is the ideal term: $`e`$ is the base of natural logarithms and the $`\rho _i^{}`$ are unimportant base units of concentration related to the definition of the standard state Smith (1990). The second term is a mean-field electrostatics term: $`\rho _z(𝐫)=_iz_i\rho _i(𝐫)`$ is the local charge density with $`z_i=\{Z,1,1\}`$ as $`i=\{m,+,\}`$, and a factor $`1/2`$ allows for double counting. The third term (correlation term) represents the excess free energy. As discussed above, only the macroion self energy $`f_m`$ is included in this term. This is computed using Debye-Hückel theory van Roij and Hansen (1997); van Roij et al. (1999); Beresford-Smith et al. (1985); not (a),
$$f_m(𝐫)=\frac{2Z^2l_\mathrm{B}k_\mathrm{B}T}{\sigma (\sigma \kappa (𝐫)+2)},$$
(2)
where $`\kappa (𝐫)`$ is a local inverse Debye screening length. This is defined in terms of an *average* local ionic strength, $`\overline{\rho }_I(𝐫)`$, through
$$\begin{array}{c}[\kappa (𝐫)]^2=8\pi l_\mathrm{B}\overline{\rho }_I(𝐫),\hfill \\ \overline{\rho }_I(𝐫)=d^3𝐫^{}w(|𝐫𝐫^{}|)\rho _I(𝐫^{})\hfill \\ \rho _I(𝐫^{})=[\rho _+(𝐫^{})+\rho _{}(𝐫^{})]/2.\hfill \end{array}$$
(3)
The ionic strength includes the counterions and salt ions, but not the macroions. In principle, allowance should be made for the macroion excluded volume, but this effect is of secondary importance and for simplicity has been omitted.
The smoothing kernel in the second of Eqs. (3) is normalised so that $`d^3𝐫w(r)=1`$. Here I use
$$w(r)=(\pi \alpha \sigma ^2)^{3/2}\mathrm{exp}[r^2/(\alpha \sigma ^2)].$$
(4)
This is an arbitrarily chosen function not (b), of range $`\alpha ^{1/2}\sigma `$. The argument below suggests that the parameter $`\alpha `$ should be of order unity and for the most part I will set $`\alpha =1`$ in the calculations. Eqs. (1)–(4) completely specify the density functional theory, and everything discussed below can be derived from them.
The decomposition into ideal, mean field, and correlation contributions is a standard approach Evans and Sluckin (1980); Sluckin (1981); dec . The approximation made for the correlation term deserves more discussion though. The only piece of physics that has been incorporated is the macroion self energy. This has a non-trivial dependence on the local ionic strength since each macroion polarises the surrounding electrolyte and becomes surrounded by a ‘double layer’. This dependence causes macroions to drift towards regions of high ionic strength, as discussed already.
The physical reason for introducing a smoothing kernel is that one can derive the self energy by integrating out the small ion degrees of freedom, with the main contribution coming from variations on length scales corresponding to the structure in the double layer van Roij et al. (1999). Thus only variations in ionic strength on length scales $`\sigma `$ should be included in the model. The smoothing kernel is a device for achieving this. This argument also motivates the choice for $`\alpha `$ in Eq. (4).
In section V below, it is found that the theory is not well behaved if one uses a ‘point model’ where the dependence is on the ionic strength at, say, the centre of the macroion (the first of Eqs. (3) with $`\overline{\rho }_I`$ replaced by $`\rho _I`$). This provides a second technical reason to make the self energy depend on a smeared ionic strength.
The potential energy of a small ion at the surface of the macroion, in units of $`k_\mathrm{B}T`$, is $`\pm Zl_\mathrm{B}/\sigma `$. Eq. (2) uses the Debye-Hückel expression for the self energy, which assumes $`Zl_\mathrm{B}/\sigma 1`$. The expression becomes increasingly inaccurate for $`Zl_\mathrm{B}/\sigma 1`$, and its use has been the subject of strong criticism as discussed above. Since the interesting effects are found only at larger values of $`Zl_\mathrm{B}/\sigma `$, one should interpret the quantitative results with caution.
## III Bulk phase behaviour
In this section, I shall consider the bulk phase behaviour predicted by the free energy of Eqs. (1)–(4). This is a homogeneous situation in which the density variables lose their spatial dependence. In this limit, one can prove that the mean field term should be replaced by a condition of bulk charge neutrality, $`\rho _z=_iz_i\rho _i=0`$ Warren (1997); Tamashiro and Schiessel (2003).
The required charge neutrality condition can be imposed in two ways. The first route is to add a term $`\psi k_\mathrm{B}T_iz_i\rho _i`$ to the free energy, where $`\psi k_\mathrm{B}T`$ is a Lagrange multiplier. This approach has the advantage of making a close connection to the density functional theory. Taking this approach, the free energy becomes
$$\frac{F}{Vk_\mathrm{B}T}=\underset{i}{}\rho _i(\mathrm{ln}\frac{\rho _i}{e\rho _i^{}}+z_i\psi )+\frac{2Z^2l_\mathrm{B}\rho _m}{\sigma (\sigma \kappa +2)}$$
(5)
where $`V`$ is the system volume and $`\kappa ^2=4\pi l_\mathrm{B}(\rho _++\rho _{})`$. The distinction between the smoothed and unsmoothed ionic strength disappears in the homogeneous limit. In this approach the $`\rho _i`$ are treated as three independent density variables. At the end of any calculations, $`\psi `$ is adjusted to get $`_iz_i\rho _i=0`$. The value of $`\psi `$ depends on the state point under consideration.
The second way to enforce charge neutrality is to eliminate one of the density variables. Since this is numerically quite convenient, it is the approach that shall be adopted in the rest of this section. At this point one can recognise that the coions come from added salt and write $`\rho _{}=Z\rho _m+\rho _s`$ and $`\rho _+=\rho _s`$, where $`\rho _s`$ is the added salt concentration. The free energy is given by Eq. (5) but with $`\psi =0`$, and $`\rho _\pm `$ substituted by the above expressions. There are now only two independent density variables and the phase behaviour can be represented in the $`(\rho _m,\rho _s)`$ plane.
I now discuss the phase behaviour predicted by this free energy. Firstly, in the absence of salt some additional simplifications can be made. In the limit $`\rho _s0`$, the free energy can be written in a dimensionless form as
$$\frac{\pi \sigma ^3F}{6ZVk_\mathrm{B}T}=\varphi \mathrm{ln}\varphi +\frac{2\varphi Zl_\mathrm{B}/\sigma }{(24\varphi Zl_\mathrm{B}/\sigma )^{1/2}+2}$$
(6)
where $`\varphi `$ is the macroion volume fraction. To get to this point, I have assumed that $`Z1`$ and hidden some constants and terms strictly proportional to $`\rho _m`$ since they do not affect the phase behaviour.
Eq. (6) predicts the dependence on $`\sigma /l_\mathrm{B}`$ and $`Z`$ is through the single combination $`Zl_\mathrm{B}/\sigma `$ (there is no reason to suppose that this should be the case in a more accurate theory). This is the same parameter that quantifies the accuracy of the Debye-Hückel linearisation approximation. The inverse of this, $`\sigma /(Zl_\mathrm{B})`$, is proportional to the dimensionless temperature discussed above. Fig. 1(a) shows the universal phase behaviour predicted by Eq. (6) as a function of the macroion volume fraction and $`\sigma /(Zl_\mathrm{B})`$. At small enough values of $`\sigma /(Zl_\mathrm{B})`$, a two phase region is encountered in the phase diagram. The two phase region corresponds to phase coexistence between macroion rich and macroion poor phases. The identities of these phases merge at a critical point located at $`\varphi 9.18\times 10^3`$ and $`\sigma /(Zl_\mathrm{B})0.132`$.
One can compare this with the simulation results of Reščič and Linse for $`Z=10`$ macroions Reščič and Linse (2001). They also find a two phase region on lowering temperature, with a critical point located at $`\varphi 0.17`$ and $`\sigma /(Zl_\mathrm{B})0.077`$. Whilst the phenomenology is the same, the numerical values are somewhat different from the prediction of Eq. (6). Not unexpectedly, the present model is too crude to obtain quantitatively reliable results. An analogy can be made with the application of Debye-Hückel theory to the restricted primitive model (RPM) Fisher and Levin (1993); Fisher (1994); Luijten et al. (2002). In this case too, Debye-Hückel theory correctly suggests a region of phase separation at low temperatures but errs in terms of quantitative predictions. Interestingly, in terms of accuracy of prediction, the present theory is not much worse than symmetrised Poisson-Boltzmann theory or the mean spherical approximation González-Tovar (1999); Bhuiyan and Outhwaite (2002); not (c).
I now turn the effect of added salt, and analyse the predictions of the full free energy in Eq. (5). In general, as salt is added, the critical point in Fig. 1(a) first moves to higher dimensionless temperatures, passes through a maximum, and then starts to move to lower dimensionless temperatures again. This non-monotonic behaviour is shown in Fig. 1(b) for $`Z=10^3`$. A similar effect of added salt is seen in a number of other approaches van Roij and Hansen (1997); van Roij et al. (1999); Warren (2000); Diehl et al. (2001); Petris and Chan (2002); Dufrêche et al. (2003). In the presence of added salt, it is no longer true that the dependence on $`Z`$ and $`\sigma /l_\mathrm{B}`$ can be combined into a single parameter, however for comparison with the phase behaviour in the absence of salt, Fig. 1(b) shows the behavior as a function of $`\sigma /(Zl_\mathrm{B})`$ at this fixed value of $`Z`$.
The re-entrant behaviour means that for parameters such as those corresponding to the dashed line in Fig. 1(b), there are *two* critical points in the $`(\rho _m,\rho _s)`$ plane, and one encounters a re-entrant single phase region at low added salt. The dashed line in Fig. 1(b) is for $`Z=10^3`$, $`\sigma =100\mathrm{nm}`$ and $`l_\mathrm{B}=0.72\mathrm{nm}`$, and the corresponding phase behaviour in the $`(\rho _m,\rho _s)`$ plane is shown in Fig. 2. It is seen that the two phase region appears as a miscibility gap in this representation.
As $`\sigma /l_\mathrm{B}`$ is increased or $`Z`$ is decreased, the two critical points move towards each other and finally disappear at a double critical point, or hypercritical point Walker and Vause (1983). For example, for $`Z=10^3`$ the double critical point corresponds to the maximum of the solid line in Fig. 1(b), where $`\sigma /(Zl_\mathrm{B})0.145`$, $`\varphi 1.04\times 10^2`$ and $`\rho _s8.98\mu \mathrm{M}`$.
The bulk phase behaviour predicted by Eq. (5) thus closely resembles that predicted by various other approaches, including the theory discussed in Ref. Warren (2000). Many approaches, including the present one, do not consider the formation of ordered phases (colloidal crystals). These can arise from the strong macroion-macroion interactions. The possibility of ordered phases has been considered by van Roij and coworkers van Roij and Hansen (1997); van Roij et al. (1999); Hynninen et al. (2004) though. They find that ordered phases can appear in the vicinity of the miscibility gap in which case a richer phase behaviour can result.
## IV Interfacial properties
A major use of the density functional theory in the present context is to calculate the macroion and small ion density profiles through the interface between two coexisting phases, and to compute the surface tension. In order to set the problem up, it is convenient to introduce the grand potential Eva
$$\mathrm{\Omega }=Fd^3𝐫\underset{i=m,\pm }{}\mu _i\rho _i(𝐫)$$
(7)
where $`\mu _i`$ are the chemical potentials of the three species, and $`F`$ is defined in Eqs. (1)–(4). At this point it is also convenient to rewrite the mean field term in Eq. (1). Define a dimensionless electrostatic potential
$$\psi (𝐫)=l_\mathrm{B}d^3𝐫^{}\frac{\rho _z(𝐫^{})}{|𝐫𝐫^{}|}$$
(8)
so that the mean field term in Eq. (1) can be written
$$\frac{l_\mathrm{B}}{2}d^3𝐫d^3𝐫^{}\frac{\rho _z(𝐫)\rho _z(𝐫^{})}{|𝐫𝐫^{}|}=\frac{1}{2}d^3𝐫\psi (𝐫)\rho _z(𝐫).$$
(9)
By direct substitution, one verifies that the potential defined by Eq. (8) solves the Poisson equation
$$^2\psi +4\pi l_\mathrm{B}\rho _z=0.$$
(10)
Using this and Green’s first identity Hinchey (1980), the mean field term now becomes
$$\frac{1}{2}d^3𝐫\psi (𝐫)\rho _z(𝐫)=\frac{1}{8\pi l_\mathrm{B}}d^3𝐫|\psi |^2.$$
(11)
This is recognised as the electric field energy since $`\psi `$ is essentially the electric field strength. One can now define a grand potential density $`\omega (𝐫)`$ such that $`\mathrm{\Omega }=d^3𝐫\omega (𝐫)`$ and
$$\omega =\underset{i}{}\rho _i\left(k_\mathrm{B}T\mathrm{ln}\frac{\rho _i}{e\rho _i^{}}\mu _i\right)+\frac{k_\mathrm{B}T}{8\pi l_\mathrm{B}}|\psi |^2+f_m\rho _m$$
(12)
where the explicit dependence on the spatial co-ordinate has been suppressed. For a homogeneous system, $`\omega =p`$ where $`p`$ is the pressure.
Setting $`\delta \mathrm{\Omega }/\delta \rho _i(𝐫)=0`$ and using Eq. (8) gives
$$\begin{array}{c}\frac{\mu _i}{k_\mathrm{B}T}=\mathrm{ln}\frac{\rho _i(𝐫)}{\rho _i^{}}+z_i\psi (𝐫)\hfill \\ +\frac{\delta }{\delta \rho _i(𝐫)}\left(\frac{d^3𝐫^{}\rho _m(𝐫^{})f_m(𝐫^{})}{k_\mathrm{B}T}\right).\hfill \end{array}$$
(13)
In principle, these non-linear integral equations can be solved to find the ion density profiles. Here a variational approximation has been adopted in which $`\mathrm{\Omega }`$ is minimised with respect to parameters in trial functions which specify the ion density profiles. More details of the numerical approach are given in Appendix B.
I now suppose that all the variation occurs in one direction $`x`$ normal to the interface. At large distances from the interface, $`x\pm \mathrm{}`$, the number densities approach those corresponding to the coexisting bulk phases. The grand potential density approaches a constant value $`\omega (\pm \mathrm{})`$ equal to (minus) the pressure, and therefore the same in coexisting phases. The surface tension $`\gamma `$ can therefore be identified as the excess grand potential per unit area
$$\gamma =_{\mathrm{}}^{\mathrm{}}𝑑x[\omega \omega (\pm \mathrm{})].$$
(14)
The chemical potentials derived from Eq. (5) are
$$\frac{\mu _i}{k_\mathrm{B}T}=\mathrm{ln}\frac{\rho _i}{\rho _i^{}}+z_i\psi +\frac{}{\rho _i}\left(\frac{2Z^2l_\mathrm{B}\rho _m}{\sigma (\sigma \kappa +2)}\right).$$
(15)
Comparison with Eq. (13) shows that $`\psi `$ in this expression is simply the limiting value of $`\psi (𝐫)`$ in the case of a homogeneous system not (d). For the interface problem, one has two limiting values, $`\psi (\pm \mathrm{})`$. The difference $`\mathrm{\Delta }\psi =\psi (\mathrm{})\psi (\mathrm{})`$ arises because of the electrical structure at the interface. It is a liquid-liquid junction potential analogous to the Donnan potential that appears across a semi-permeable membrane Don . Since $`\psi `$ in Eq. (15) is determined by the bulk densities, the difference $`\mathrm{\Delta }\psi `$ can be calculated without having to solve for the interface structure. In fact, because of the symmetric way that $`\rho _\pm `$ enters into the excess free energy, a simple expression obtains,
$$\mathrm{\Delta }\psi =\frac{1}{2}\mathrm{ln}\left(\frac{\rho _{}(\mathrm{})}{\rho _+(\mathrm{})}\frac{\rho _+(\mathrm{})}{\rho _{}(\mathrm{})}\right).$$
(16)
This method of calculating the junction potential was used in Ref. Warren (2000).
One question remains: what should be used for the chemical potentials in these calculations? The simplest answer is to compute the chemical potentials from Eq. (15), setting $`\psi =0`$ and using the bulk densities corresponding to either one of the coexisting phases. This works because global charge neutrality means Eq. (14) for the surface tension is unaffected by the value of $`\psi `$ in Eq. (15). Hence we are free to set $`\psi =0`$ in either of the coexisting phases.
I now turn to the results. Fig. 3 shows representative density profiles for the macroion and small ions through the interface between the coexisting phases, corresponding to the highlighted tie line in Fig. 2. The profiles interpolate smoothly between the coexisting bulk densities. Fig. 4 shows the detailed electrical structure at the interface. The upper plot shows that the charge density $`\rho _z=Z\rho _M+\rho _+\rho _{}`$ has a dipolar structure. Correspondingly there is a localised electric field, shown in the middle plot, and a smooth jump of $`\mathrm{\Delta }\psi 20.5\mathrm{mV}`$ in the electrostatic potential, shown in the lower plot. This is the junction potential which can also be calculated directly from the coexisting bulk densities as in Eq. (16). This electrical structure is in accord with general expecatations for charged systems Sluckin (1981); Nab .
Fig. 5 shows the grand potential density and the electrostatic component thereof—the second term of Eq. (12)—as a function of distance through the interface. For this particular case the area gives $`\gamma 0.727\times (k_\mathrm{B}T/\sigma ^2)`$. The order of magnitude of this should not come as a surprise since $`\sigma `$ and $`k_\mathrm{B}T`$ are the only relevant length and energy scales in the problem. Inserting actual values, $`\gamma 0.3\mu \mathrm{N}\mathrm{m}^1`$, which is typical for for soft matter interfaces sur
Fig. 6 shows how the surface tension and interface width vary as one approaches the upper critical point in Fig. 2. The width $`d`$ is defined operationally as $`d^2=x^2x^2`$, where $`\mathrm{}=_{\mathrm{}}^{\mathrm{}}(\mathrm{})p(x)𝑑x/_{\mathrm{}}^{\mathrm{}}p(x)𝑑x`$, with $`p(x)=|\omega (x)\omega (\pm \mathrm{})|^2`$. These results are obtained by repeating the calculations underlying Figs. 35 for a sequence of tie lines approaching the critical point. They are reported as a function of the distance from the critical point, expressed in terms of a normalised salt chemical potential. Fig. 6(b) also shows the correlation lengths $`\xi _\pm `$ in the coexisting phases determined from the exponential decay of the density profiles into the bulk phases (see Appendix B). As the critical point is approached, these approach each other, and diverge in the same way as the interface width. Fig. 6 reveals that the surface tension and length scales are in accord with expected scaling behaviour for a mean-field theory Rowlinson and Widom (1989).
What happens at the lower critical point in Fig. 2 though? The next section shows that this is a non-trivial question with perhaps an unexpected answer. In the calculations in the current section, I have assumed that the interface profiles smoothly interpolate between the coexisting phases. Indeed, this is the basis of the numerical method detailed in Appendix B. However, such an approach rules out the possibility of oscillatory behaviour in the density profiles (or to be precise, the numerical methodology is inappropriate for this scenario). At lower salt concentrations though, one can enter a region where oscillatory behaviour is expected. These considerations are made mathematically precise in the next section.
## V Structure factors
The structure factors in a homogeneous system can be determined from a density functional theory (DFT) by functional differentiation Eva . Where accurate structure factors are already known, typically from a combination of simulation and integral equation approaches, this can be used to constrain the DFT. In the present case for example, one could try to constrain $`w(r)`$ in Eq. (4). However accurate structure factors are not known for this problem, and furthermore the DFT has been constructed to include only the macroion self energy. Thus it does not make sense to constrain the DFT and the present section simply reports the structure factors that are predicted from the theory as given in Eqs. (1)–(4).
The structure factor matrix is Hansen and McDonald (1976); March and Tosi (1976)
$$\stackrel{~}{S}_{ij}(q)=\rho _i\delta _{ij}+\rho _i\rho _j\stackrel{~}{h}_{ij}(q)$$
(17)
where $`i`$ and $`j`$ run over $`\{m,+,\}`$ and $`\stackrel{~}{h}_{ij}(q)=d^3𝐫e^{i𝐪𝐫}h_{ij}(r)`$ is the Fourier transform of the pair correlation functions $`h_{ij}(r)=g_{ij}(r)1`$. Reciprocal space quantities will be denoted by a tilde. The bulk densities $`\rho _i`$ are constants, fixed by the choice of state point. Deviations away from these will be denoted by $`\mathrm{\Delta }\rho _i`$. Eq. (17) uses the normalisation $`\stackrel{~}{S}_{ij}(q)\rho _i\delta _{ij}`$ as $`q\mathrm{}`$, which simplifies some of the expressions below March and Tosi (1976).
To obtain the structure factor matrix, start by defining the real-space function
$$S_{ij}^1(|𝐫𝐫^{}|)=\frac{1}{k_\mathrm{B}T}\left(\frac{\delta ^2F}{\delta \rho _i(𝐫)\delta \rho _j(𝐫^{})}\right)_{\rho _i(𝐫)\rho _i}$$
(18)
where $`F`$ is the full free energy. The limit of a homogeneous system is taken after the functional differentiation step so that $`S_{ij}^1`$ only depends on $`|𝐫𝐫^{}|`$ as indicated. Transforming to reciprocal space, one can show that
$$\stackrel{~}{S}_{ij}^1(q)=d^3𝐫e^{i𝐪𝐫}S_{ij}^1(r)$$
(19)
is simply the matrix inverse of $`\stackrel{~}{S}_{ij}`$,
$$\underset{j}{}\stackrel{~}{S}_{ij}\stackrel{~}{S}_{jk}^1=\delta _{ik}.$$
(20)
These results follow by combining the Ornstein-Zernike relation for a multicomponent mixture in reciprocal space, $`\stackrel{~}{h}_{ij}=\stackrel{~}{c}_{ij}+_k\rho _k\stackrel{~}{c}_{ik}\stackrel{~}{h}_{jk}`$ where $`c_{ij}`$ are the direct correlation functions Hansen and McDonald (1976), with the DFT result that $`c_{ij}=(1/k_\mathrm{B}T)\delta ^2F_{\mathrm{ex}}/\delta \rho _i\delta \rho _j`$ where $`F_{\mathrm{ex}}`$ is the excess free energy Eva .
The route to the structure factors offered by Eqs. (18)–(20) is based on ‘classical’ arguments Hansen and McDonald (1976). One can also make the connection via field theoretical methods. Expanding the free energy functional to second order gives
$$\frac{\mathrm{\Delta }F}{k_\mathrm{B}T}=\frac{1}{2}d^3𝐫d^3𝐫^{}\underset{ij}{}\mathrm{\Delta }\rho _i(𝐫)\mathrm{\Delta }\rho _j(𝐫^{})S_{ij}^1(|𝐫𝐫^{}|),$$
(21)
where $`S_{ij}^1`$ is defined by Eq. (18). It follows that Doi and Edwards (1986)
$$\mathrm{\Delta }\rho _i(𝐫)\mathrm{\Delta }\rho _j(𝐫^{})=S_{ij}(|𝐫𝐫^{}|)$$
(22)
where $`S_{ij}(r)=d^3𝐪/(2\pi )^3e^{i𝐪𝐫}\stackrel{~}{S}_{ij}(q)`$ is the structure factor matrix expressed as a real space quantity. Although care has to be taken at the point $`𝐫=𝐫^{}`$, one can easily show that the density-density correlation function on the left hand side of Eq. (22) is the same as the Fourier transform of the right hand side of Eq. (17).
The Stillinger-Lovett moment conditions constrain the behaviour of the structure factors in reciprocal space in a particularly clear manner Stillinger and Lovett (1968); Evans and Sluckin (1980); Martin (1988); Stafiej and Badiali (1997); Lee and Fisher (1997). Firstly, the zeroth-moment conditions express perfect screening and are $`d^3𝐫_iz_i\rho _ig_{ij}(r)=z_j`$ for $`j=\{m,+,\}`$. Using charge neutrality and assuming the structure factors are regular at $`q=0`$, one can easily show that this implies
$$_iz_i\stackrel{~}{S}_{ij}(𝐪)=O(q^2).$$
(23)
The second-moment condition is $`d^3𝐫r^2_{ij}z_iz_j\rho _i\rho _jg_{ij}(r)=3/(2\pi l_\mathrm{B})`$. This constrains the long wavelength behaviour of the charge-charge structure factor,
$$_{ij}z_iz_j\stackrel{~}{S}_{ij}(𝐪)=\frac{q^2}{(4\pi l_\mathrm{B})}+O(q^4).$$
(24)
In real space, this means that $`\mathrm{\Delta }\rho _z(𝐫)\mathrm{\Delta }\rho _z(𝐫^{})l_\mathrm{B}/|𝐫𝐫^{}|`$ for $`|𝐫𝐫^{}|\mathrm{}`$. Thus charge density fluctuations vanish with the Coulomb law at large distances, corresponding to the fact that the electrostatic energy dominates in the free energy for long-wavelength density fluctuations unless they happen to be charge-neutral not (e).
I now apply the formalism of Eqs. (18)–(20) to the present DFT defined in Eqs. (1)–(4). The result for the inverse structure factor matrix in reciprocal space can be written as
$$\stackrel{~}{S}_{ij}^1=\stackrel{~}{T}_{ij}^1+\frac{4\pi l_\mathrm{B}z_iz_j}{q^2}$$
(25)
where first term comes from the ideal and correlation contributions to the free energy and the second term from the mean field electrostatics. The first term is in detail
$$\begin{array}{c}\stackrel{~}{T}_{ij}^1=\delta _{ij}/\rho _i+\rho _mZ^2\pi ^2l_\mathrm{B}^3\sigma ^3h_1(\sigma \kappa ,\sigma q)\mathrm{\Delta }_{ij}^{}\hfill \\ Z^2\pi l_\mathrm{B}^2\sigma h_2(\sigma \kappa ,\sigma q)\mathrm{\Delta }_{ij}^{\prime \prime }.\hfill \end{array}$$
(26)
where the functions $`h_{1,2}(x=\sigma \kappa ,y=\sigma q)`$ are
$$h_1=\frac{8e^{\alpha y^2/2}(2+3x)}{(x^3(x+2)^3)},h_2=\frac{4e^{\alpha y^2/4}}{(x(x+2)^2)},$$
(27)
and the matrices are
$$\begin{array}{cc}\mathrm{\Delta }_{mm}^{}=\mathrm{\Delta }_{m\pm }^{}=0,\hfill & \mathrm{\Delta }_{\pm \pm }^{}=1,\hfill \\ \mathrm{\Delta }_{mm}^{\prime \prime }=\mathrm{\Delta }_{\pm \pm }^{\prime \prime }=0,\hfill & \mathrm{\Delta }_{m\pm }^{\prime \prime }=1.\hfill \end{array}$$
(28)
The $`y`$-dependence ($`y=\sigma q`$) in Eq. (27) arises from the Fourier transform of the weight function of Eq. (4). Note that the point model alluded to in section II corresponds to the limit $`\alpha 0`$ in Eqs. (27). In this limit, the theory becomes ill-defined since $`\stackrel{~}{S}_{ij}(q)`$ does not have the correct limiting behaviour as $`q\mathrm{}`$. This was the original technical reason for introducing the smoothing kernel.
For any given state point and value of $`q`$, Eqs. (25)–(28) define $`\stackrel{~}{S}_{ij}^1`$ which can be inverted numerically to find all components of the structure factor matrix. A partial solution can be obtained analytically in terms of the subsidiary matrix $`\stackrel{~}{T}_{ij}`$,
$$\stackrel{~}{S}_{ij}=\stackrel{~}{T}_{ij}\frac{4\pi l_\mathrm{B}_{kl}z_kz_l\stackrel{~}{T}_{ik}\stackrel{~}{T}_{jl}}{q^2+4\pi l_\mathrm{B}_{kl}z_kz_l\stackrel{~}{T}_{kl}}$$
(29)
From this one can readily prove that $`\stackrel{~}{S}_{ij}`$ exactly satisfies the Stillinger-Lovett moment conditions in Eqs. (23) and (24) above.
Another result follows from the dominance of the ideal contribution over the correlation contribution at low densities. In the limit $`\rho _i0`$ one finds $`\stackrel{~}{T}_{ij}\rho _i\delta _{ij}`$ and
$$\stackrel{~}{S}_{ij}\rho _i\delta _{ij}\frac{4\pi l_\mathrm{B}z_iz_j\rho _i\rho _j}{q^2+4\pi l_\mathrm{B}_kz_k^2\rho _k}.$$
(30)
This is in fact exactly in accordance with the Debye-Hückel limiting law at low densities. To see this, note that $`\lambda =(4\pi l_\mathrm{B}_kz_k^2\rho _k)^{1/2}`$ is the Debye screening length defined to include *all* ionic species. Thus in real space, Eqs. (17) and (30) indicate that $`h_{ij}=z_iz_j(l_\mathrm{B}/r)e^{r/\lambda }`$, in correspondence with the Debye-Hückel limiting law.
It is clear that the moment conditions and the Debye-Hückel limiting law behaviour follow from the construction of the DFT to include a mean-field contribution separately from the correlation term. This construction is in turn motivated by the expected behaviour of the direct correlation functions $`c_{ij}(r)`$ at $`r\mathrm{}`$, as Evans and Sluckin have described Evans and Sluckin (1980). The form of the correlation term is unimportant, so long as it is regular both at $`q0`$ and $`\rho _i0`$.
For the remaining part, I now focus on the macroion structure factor $`\stackrel{~}{S}_{mm}`$. Note that the theory includes the macroion-macroion electrostatic interaction explicitly in the mean field term, and an additional indirect interaction in the correlation term. The computation of $`\stackrel{~}{S}_{mm}`$ reveals the combined effect of these macroion *interactions* on the macroion *correlations*.
Typically $`\stackrel{~}{S}_{mm}`$ has a ‘hole’ in reciprocal space for $`q\sigma 1`$. This corresponds to the macroion electrostatic repulsions. Within the correlation hole though, there is additional structure. This becomes particularly important in the vicinity of the phase separation region. Two kinds of behaviour are possible: at higher salt concentrations $`\stackrel{~}{S}_{mm}`$ rises to a maximum as $`q0`$, or at lower salt concentrations $`\stackrel{~}{S}_{mm}`$ acquires a peak at some $`q^{}>0`$. In the phase diagram, the two alternatives are separated by a (macroion) ‘Lifshitz line’ Archer et al. (2002), defined to be the locus of points for which $`\stackrel{~}{S}_{mm}/(q^2)|_{q=0}=0`$. Fig. 7(a) shows the two behaviours for a pair of typical state points above and below the Lifshitz line, and Fig. 7(b) shows the Lifshitz line superimposed on the bulk phase behaviour.
Also shown in Fig. 7(b) is the spinodal line computed from the bulk free energy in Eq. (5) of section III. One can check that $`\stackrel{~}{S}_{mm}(q=0)`$ diverges on this spinodal line; in fact all the $`q=0`$ components of the structure factor matrix diverge because the determinant of $`\stackrel{~}{S}_{ij}^1`$ vanishes. For salt concentrations above the Lifshitz line, this divergence at $`q=0`$ can be accommodated within the general behaviour of the structure factor. Of course, state points within the binodal are metastable so the divergence is strictly only visible as the upper critical point is approached. The fact that the structure factors diverge on the spinodal line is no coincidence, since thermodynamic consistency by the compressibility route is assured for a DFT not (f).
What happens at salt concentrations below the Lifshitz line? Here, the peak in $`S_{mm}`$ at $`q^{}>0`$ is found to diverge *before* the bulk spinodal line is reached. The shaded area in Fig. 7(b) shows the region where this occurs. A divergence at a non-zero wavevector is indicative of microphase separation mic . In this case one would expect a charge-density-wave (CDW) phase to appear Nabutovskii et al. (1980, 1985). The shaded region extends below the binodal for bulk phase separation, so the CDW phase should be observable in this part of the phase diagram. In fact the CDW phase will be found whenever the lower critical point lies below the Lifshitz line. The general idea that a critical point in a charged system can be replaced by a CDW phase was advanced by Nabutovskii, Nemov and Peisakhovich Nabutovskii et al. (1980); not (g).
The location of the Lifshitz line depends on the parameter $`\alpha `$ which sets the range of the smoothing kernel $`w(r)`$ in Eq. (4). If $`\alpha 0.40`$ the Lifshitz line moves upwards past the upper critical point, which would then be expected to be replaced by a CDW phase too. On the other hand if $`\alpha 3.6`$, the Lifshitz line moves downwards past the lower critical point. These critical values of $`\alpha `$ only depend on the coefficient of $`q^2`$ in the expansion of the Fourier transform of $`w(r)`$ about $`q=0`$.
The Lifshitz line discussed here pertains to the macroion structure factor. Although slightly different Lifshitz lines are expected for each component of the structure factor matrix, the locus of state points where the peak diverges (either on the spinodal or on the boundary of the CDW phase) should be the same for all components.
Whilst the Lifshitz line line marks an obvious change in the behaviour of $`\stackrel{~}{S}_{mm}`$, the cross-over from monotonic to damped oscillatory asymptotic decay of the correlation functions $`h_{ij}(r)`$ is determined by Kirkwood or Fisher-Widom lines in the phase diagram KFW ; Evans et al. (1994); Leote de Carvalho and Evans (1994). The difference between these is rather subtle Leote de Carvalho and Evans (1994); not (h), and one might loosely cover both possibilities by the phrase ‘Kirkwood-Fisher-Widom’ (KFW) line. The importance of the KFW line lies in the fact that it also governs the asymptotic decay of the interface density profiles, which behave in the same way as $`h_{ij}`$ Evans et al. (1994). Thus the calculations reported in section IV above, which assume that there is no oscillatory behaviour in the density profiles, requires as a necessary minimum that the coexisting bulk densities both lie above the KFW line. The location of the KFW line is governed by the poles of $`\stackrel{~}{S}_{ij}(q)`$ in the complex $`q`$ plane, which are either purely imaginary or occur as complex conjugate pairs, and are the same for all components of $`\stackrel{~}{S}_{ij}`$ Evans et al. (1994). If the pole nearest the real $`q`$-axis is purely imaginary, then monotonic decay is expected; conversely if a pair of complex conjugate poles is nearest the real $`q`$-axis, then damped oscillatory decay is expected Leote de Carvalho and Evans (1994). Determination of the KFW line is a hard numerical problem and has not been attempted for the present DFT. However the presence of a peak in $`\stackrel{~}{S}_{mm}(q)`$ on the real $`q`$-axis at $`q=0`$, or at $`q^{}>0`$, ought to be indicative of whether the pole nearest the real $`q`$-axis is, or is not, purely imaginary. Thus the Lifshitz line should serve as a guide to the location of the KFW line. In section IV therefore, care was taken to make sure that the coexisting bulk densities lie well above the Lifshitz line.
## VI Discussion
The paper presents a density functional theory (DFT) for a macroion suspension. The excess free energy corresponds to the macroion self energy evaluated using Debye-Hückel theory. These approximations render theory tractable without losing the basic phenomenology which resembles that of other studies. The advantage of a DFT is that one can compute the interface structure and surface tension between coexisting phases. The results are in accord with expectations from previous work Warren (2000). In particular, the electrical structure of the interface gives rise to a junction potential analogous to the Donnan potential across a semi-permeable membrane. This arises from an electric dipole moment density (per unit area of interface), which appears because charge neutrality is locally violated in the vicinity of the interface. The surface tension is found to be of the order $`k_\mathrm{B}T/\sigma ^2`$.
Structure factors can be computed from the DFT. These are found to obey the Stillinger-Lovett moment conditions, although this is not a stringent test of the theory. The structure factors reveal an interesting phenomenon, namely that oscillatory behaviour can appear in the (direct) correlation functions, particularly at low ionic strength. Indeed there may be regions of microphase separation in the vicinity of the critical points, corresponding to the appearance of a charge-density-wave (CDW) phases. This phenomenon is peculiar to asymmetric charged systems Nabutovskii et al. (1980), and is strictly absent in symmetric systems such as the restricted primitive model. In this respect, the possibility of CDW phases is correlated with the appearance of the junction potential, which is also strictly absent in symmetric systems War . Given the approximate nature of the DFT, only certain aspects of the present analysis might be expected to survive in a full treatment. One of these is an upturn in macroion structure factor at small $`q`$, even in the absence of a true miscibility gap. This would reflect an increased osmotic compressibility in this region of the phase diagram. Another expectation is the possible appearance of the CDW phases, although it might be difficult to disentangle these from the ordered (crystal) phases that are expected for a macroion suspension at sufficiently strong electrostatic coupling.
The macroion self energy depends on the local ionic strength, but on both physical and technical grounds it is found necessary to introduce the notion of smoothing or smearing—the dependency should be on the ionic strength averaged over the vicinity of the macroion. Here a completely phenomenological approach has been taken to construct the details of the DFT. Other choices could be made, or indeed more rigor could be introduced, such as additional requirements for internal consistency not (i). Tests indicate though that the general phenomenology (electrical structure of interface, gross behaviour of structure factors) is found to be insensitive to the details of the model at this point.
###### Acknowledgements.
I thank R. Evans and A. S. Ferrante for useful discussions.
## Appendix A Correction to Ref. Warren (2000)
Chan Chan (2001) has remarked that an excluded volume contribution was omitted in the theory of Ref. Warren (2000). This appendix describes the missing term. The error occurs in going from Eq. (3) to Eq. (7) of Ref. Warren (2000) where the omitted contribution arises from the fact that $`h_{m\pm }(r)=g_{m\pm }(r)1=1`$ for $`r<\sigma /2`$. In terms of the microion-macroion interaction energy, $`E_{\mathrm{ms}}/(Vk_\mathrm{B}T)`$, the omitted contribution is
$$\begin{array}{c}\rho _m_{|𝐫|<\sigma /2}d^3𝐫\frac{Zl_\mathrm{B}}{r}[\rho _+h_{m+}(r)\rho _{}h_m(r)]\hfill \\ =\rho _m(\rho _+\rho _{})_0^{\sigma /2}4\pi r^2𝑑r\frac{Zl_\mathrm{B}}{r}\hfill \\ =+\frac{\pi Z^2l_\mathrm{B}\rho _m^2\sigma ^2}{2}\text{(using }\rho _+\rho _{}=Z\rho _m\text{).}\hfill \end{array}$$
(31)
This contribution is a positive, increasing function of $`\rho _m`$, and has the tendency to stabilise the system against phase separation (because it is an athermal excluded volume term, it passes unscathed through the thermodynamic integration step needed to calculate the contribution to the free energy). If the calculations of Ref. Warren (2000) are repeated with this contribution included, it is found that the basic phenomenology is still the same, except that the miscibility gap in the $`(\rho _m,\rho _s)`$ plane does not appear until somewhat larger values of $`Zl_\mathrm{B}/\sigma `$. Fig. 8 shows the new results in comparison with those reported in Table II of Ref. Warren (2000). The new calculation indicates that phase separation is observed in an even narrower window of parameter space for which the Debye-Hückel linearisation approximation might be admissible than was found in the earlier work. This can be taken to indicate that the self-energy mechanism may not be sufficiently powerful to drive phase separation by itself, as discussed in the introduction.
## Appendix B Numerical approach
The task is to find density profiles $`\rho _i(x)`$ which minimise the grand potential in Eq. (7). The most accurate method is to solve the integral equations for the profiles in Eq. (13). However, this is hard. An alternative is to adopt a variational approach in which $`\mathrm{\Omega }`$, or $`\gamma `$ in practice, is minimised with respect to parameters in trial functions which specify the density profiles Smi . This is the approach that has been taken here.
The ion density profiles have to satisfy a sum rule since the potential difference $`\mathrm{\Delta }\psi =\psi (\mathrm{})\psi (\mathrm{})`$ is fixed by the coexisting bulk densities as described in section IV. One can replace one of the ion density profiles by $`\psi (x)`$ to ensure this sum rule is automatically satisfied. In the present case, a choice was made to use the set $`\{\rho _m,\rho _+,\psi \}`$ as a basis with $`\rho _{}`$ derived analytically from the Poisson equation, $`\rho _{}=Z\rho _m+\rho _{}(d^2\psi /dx^2)/(4\pi l_\mathrm{B})`$. The first integral of the Poisson equation shows that one can additionally ensure global charge neutrality by making sure that $`d\psi /dx0`$ as $`|x|\mathrm{}`$. Once the $`\rho _i`$ are known, the average ionic strength $`\overline{\rho }_I`$ and the surface tension $`\gamma `$ are determined numerically by quadratures.
To represent the basis set $`\{\rho _m,\rho _+,\psi \}`$, three copies of the function
$$\begin{array}{c}f(x;\xi _\pm ,\{a\})=\frac{a_{}e^{x/\xi _+}a_+e^{x/\xi _{}}}{a_{}a_++a_{}e^{x/\xi _+}+a_+e^{x/\xi _{}}}\hfill \\ +_{r=1}^Na_rH_r(x/\xi )\hfill \end{array}$$
(32)
are introduced. In this, the $`H_r`$ are Hermite functions, with $`\xi =2/(1/\xi _{}+1/\xi _+)`$ used to scale the argument. Each copy of $`f`$ is parametrised by the correlation lengths $`\xi _\pm `$ and amplitude set $`\{a\}`$, and has the properties that $`f\pm (1a_\pm e^{x/\xi _\pm })`$ as $`x\pm \mathrm{}`$. One copy of $`f`$ is assigned to each member of $`\{\rho _m,\rho _+,\psi \}`$, and is scaled and shifted to match the limiting values at $`|x|\mathrm{}`$, for example $`\rho _m=\rho _m(\mathrm{})(1f)/2+\rho _m(\mathrm{})(1+f)/2`$ (for the electrostatic potential, one can set $`\psi (\mathrm{})=0`$ and $`\psi (\mathrm{})=\mathrm{\Delta }\psi `$). The three copies of $`f`$ have different amplitude sets $`\{a\}`$ but share common values for $`\xi _\pm `$ since the asymptotic decay of the density profiles into the bulk phases is expected to be governed by a bulk correlation length—it is these values of $`\xi _\pm `$ that are reported in Fig. 6(b). A finite set of $`N`$ Hermite functions has been included in each copy of $`f`$ to allow for an arbitrary structure at the interface. In practice the minimisation problem is well behaved only if the density profiles smoothly interpolate between the bulk values, for which case typically $`N=3`$–6 Hermite functions are needed to achieve convergence in $`\gamma `$ to an accuracy of the order 1%. At this point, the interface problem has been reduced to a multivariate minimisation over the three copies of the amplitude set $`\{a\}`$ plus the correlation lengths $`\xi _\pm `$. Numerical minimisation of $`\gamma `$ with respect to these parameters is then undertaken by standard methods Press et al. (1989). |
warning/0506/math0506447.html | ar5iv | text | # Unstable products of smooth curves
## 1. Introduction
A central problem in Kähler geometry is finding necessary and sufficient conditions for a Kähler class on a complex manifold $`X`$ to admit a constant scalar curvature Kähler (cscK) metric. It is well known that when the first Chern class of $`X`$ is negative, there exists a Kähler-Einstein (and hence cscK) metric in $`c_1(X)`$ . Thus by a deformation argument due to LeBrun-Simanca there exists a cscK metric in each class in an open neighborhood around $`c_1(X)`$.
It has been an element of folklore that if $`c_1(X)`$ is negative then every Kähler class on $`X`$ admits a cscK metric. We prove in this paper that this is not the case by showing that this fails on the product of certain non-generic smooth curves.
The cscK problem is related to the stability of $`X`$. A conjecture of Yau states that a rational Kähler class $`\mathrm{\Omega }`$ should admit a cscK metric if and only if the pair $`(X,\mathrm{\Omega })`$ is “stable” in the sense of geometric invariant theory . The precise definition of stability, called K-stability, was introduced by Tian and expanded by Donaldson . One direction of the conjecture has essentially be proved: the existence of a cscK metric implies stability (see (2.1)).
To show that certain Kähler classes do not admit cscK metrics we shall use an obstruction for K-stability (and hence for cscK metrics) of Thomas and the author called slope stability . The main results of this paper are the following:
Theorem 3.7. For $`g5`$ there exist smooth curves $`C`$ of genus $`g`$ such that $`X=C\times C`$ is not slope semistable with respect to certain polarisations. Thus there are Kähler classes on $`X`$ that do not admit Kähler metrics of constant scalar curvature.
Moreover this gives the first example of a manifold with negative first Chern class whose Hilbert and Chow points are unstable in the sense of geometric invariant theory.
Theorem 3.8. For $`g5`$ there exist smooth curves $`C`$ of genus $`g`$ such that $`C\times C`$ is not asymptotically Hilbert semistable (resp. not asymptotically Chow semistable) with respect to certain polarisations.
More specifically these results holds if $`C`$ admits a simple branched cover to $`^1`$ of degree $`d`$ with $`2d1<\sqrt{g}`$ (3.3). Such curves are non-generic, since a general curve of genus $`g`$ admits a branched cover to $`^1`$ of degree
$$d_0=\left[\frac{g+1}{2}\right]+1$$
and none of degree $`d<d_0`$. ( p. 261).
Given a simple branched cover $`\pi :C^1`$ of degree $`d`$, consider the fibre product $`C\times _\pi CX`$. This contains the diagonal $`\mathrm{\Delta }`$ and we let $`Z`$ be the residual divisor (i.e. $`Z+\mathrm{\Delta }=C\times _\pi C`$). We will show that if $`2d1<\sqrt{g}`$ then $`Z`$ has negative self-intersection, and the “slope” of $`Z`$ as defined in Section 2 is less than that of $`X`$, which proves instability.
As suggested by Chen, such an example leads one to try and better understand the obstruction to finding cscK metrics in these classes. In particular one technique studied by Chen , Donaldson , Weinkove and Song-Weinkove is that of the J-flow. We comment in Section 4 as to what some of these results of say when applied to products of curves.
Acknowledgments I would like to thank Richard Thomas for comments on the first draft of this paper, as well as continued encouragement. I also thank Xiuxiong Chen, Ian Morrison, Sean Paul, Michael Thaddeus and Ben Weinkove for useful conversations. Lazarsfeld’s book has been extremely helpful.
Notation By a polarisation on a variety $`X`$ we mean a choice of ample divisor which we usually denote by $`L`$ and will write $`c_1(L)`$ for the first Chern class of the associated line bundle. A $``$-divisor is a formal sum of divisors with rational coefficients, and an ample $``$-divisor is one which can be written as a sum of ample divisors (again with rational coefficients). The space of $``$-divisors modulo numerical equivalence is denoted $`N^1(X)_{}`$. Abusing notation we will often not distinguish between a $``$-divisor and its class in $`N^1(X)_{}`$.
## 2. Slope stability for Varieties
The original link between K-stability and cscK metrics is due to Tian . For slope stability we require the definition of K-stability used by Donaldson (the relation between the two definitions can be found in ).
###### Theorem 2.1.
Fix a polarised manifold $`(X,L)`$. If there exists a constant scalar curvature Kähler metric in $`c_1(L)`$ then $`(X,L)`$ is K-semistable.
###### Proof.
This is proved by in . Alternatively one can use the existence of a cscK metric to show that the Mabuchi functional is bounded from below which in turn implies K-semistability . ∎
The notion of slope stability for polarised varieties was introduced in as a necessary condition for K-stability. The general idea is that a non-generic subscheme of a polarised variety $`(X,L)`$ with certain numerical properties will have a slope that is too small, and this forces $`(X,L)`$ to be unstable.
Let $`Z`$ be a subscheme of a polarised variety $`(X,L)`$ and let $`\pi :\widehat{X}X`$ be the blowup of $`X`$ along $`Z`$ with exceptional divisor $`E`$. For sufficiently small positive $`c`$ the divisor $`\pi ^{}LcE`$ is ample; so we can define the Seshadri constant of $`Z`$ as
$$ϵ(Z,L)=\text{sup}\{c:\pi ^{}LcE\text{ is ample}\}.$$
Write the Hilbert polynomial of $`L`$ as $`\chi (kL)=a_0k^n+a_1k^{n1}+\mathrm{}`$ where $`n=dimX`$. For fixed $`x`$ we define $`a_i(x)`$ by
$$\chi (k(\pi ^{}LxE))=a_0(x)k^n+a_1(x)k^{n1}+\mathrm{}\text{for all }kx.$$
As $`\chi (k\pi ^{}LrE)`$ is a polynomial in two variables of total degree at most $`n`$, we have that $`a_i(x)`$ is a polynomial and so extends to all real $`x`$. We let $`\stackrel{~}{a}_i(x)=a_ia_i(x)`$ and for $`0<cϵ(Z,L)`$ the slope of $`X`$ and $`Z`$ is defined to be ( (3.1,3.14)),
$`\mu (X,L)`$ $`=`$ $`{\displaystyle \frac{a_1}{a_0}},`$
$`\mu _c(𝒪_Z,L)`$ $`=`$ $`{\displaystyle \frac{_0^c\stackrel{~}{a}_1(x)+\frac{\stackrel{~}{a}_0^{}(x)}{2}dx}{_0^c\stackrel{~}{a}_0(x)𝑑x}},`$
which are both finite (ibid. (4.21)).
###### Definition 2.2.
We say that $`(X,L)`$ is slope semistable with respect to a subscheme $`Z`$ if
$$\mu (X,L)\mu _c(𝒪_Z,L)\text{ for all }0<cϵ(Z,L).$$
We say $`(X,L)`$ is slope semistable if it is slope semistable with respect to all subschemes.
Since the property of being slope semistable is invariant under replacing $`L`$ by some power (ibid. (3.10)), we extend the notion of slope semistability to ample $``$-divisors.
The definition of the slopes are made so that if $`(X,L)`$ is not slope semistable with respect to $`Z`$ then the degeneration to the normal cone of $`Z`$ prevents $`(X,L)`$ being K-semistable:
###### Theorem 2.3.
If $`(X,L)`$ is not slope semistable then it is not K-semistable.
###### Proof.
See (4.18); when $`X`$ and $`Z`$ are smooth this is proved in (4.2). ∎
### Slope stability for smooth surfaces
In this paper we will only consider the case that $`X`$ is a smooth surface and $`Z`$ is a curve. Then the blowup of $`X`$ along $`Z`$ is just $`X`$ itself, so
$$ϵ(Z,L)=sup\{c:LcZ\text{ is ample}\}.$$
Letting $`K`$ be the canonical divisor of $`X`$, a simple application of the Riemann-Roch theorem to calculate $`a_0(x)`$ and $`a_1(x)`$ yields ( (5.4)),
$`\mu (X,L)`$ $`=`$ $`{\displaystyle \frac{K.L}{L^2}},`$ (2.4)
$`\mu _c(𝒪_Z,L)`$ $`=`$ $`{\displaystyle \frac{3(2L.Zc(K.Z+Z^2))}{2c(3L.ZcZ^2)}}.`$
Notice that in this case $`ϵ(Z,L)`$ as well as the slopes $`\mu _c(𝒪_Z,L)`$ and $`\mu (X,L)`$ depend only on the class of $`L`$ and $`Z`$ modulo numerical equivalence. We extend the equations (2.4) to any class $`L`$ in $`N^1(X)_{}`$ which is not necessarily ample (in which case they may no longer be finite).
### Hilbert and Chow stability
Slope stability also gives an obstruction to the classical notions of stability for projective varieties. For $`r0`$ consider the embedding of $`X`$ via the linear series $`|rL|`$ into $`^{N(r)}`$. Up to change of change of coordinates, this determines a point $`\text{Hilb}(X,L^r)`$ in the Hilbert scheme of $`^{N(r)}`$ (resp. a point $`\text{Chow}(X,L^r)`$ in the Chow variety). The action of the automorphism group of $`^{N(r)}`$ induces a linearised action on the Hilbert scheme (resp. Chow variety), and one can apply the notions of geometric invariant theory to these spaces (see ).
###### Definition 2.5.
We say that $`(X,L)`$ is asymptotically Hilbert (resp. Chow) semistable if for $`r0`$ the point $`\text{Hilb}(X,L^r)`$ (resp. $`\text{Chow}(X,L^r)`$) is semistable in the sense of geometric invariant theory.
###### Theorem 2.6.
If $`(X,L)`$ is not slope semistable then it is neither asymptotically Hilbert nor asymptotically Chow semistable.
###### Proof.
The follows from (2.3) as asymptotic Hilbert (resp. Chow) semistability implies K-semistability ( (4.18)). ∎
## 3. Unstable products of Curves
We start with some standard material on the ample cone of products of curves, all of which can be found in . Fix a smooth curve $`C`$ of genus $`g2`$ and let $`X=C\times C`$. If $`\pi _i`$ is the projection onto the $`i`$-th factor, and $`p`$ is a fixed point in $`C`$ then the class $`f_i`$ of the fibre $`\pi _i^1(p)`$ in $`N^1(X)_{}`$ is independent of $`p`$. The class of the canonical divisor of $`X`$ is $`K=(2g2)(f_1+f_2)`$ which is ample, so $`X`$ has negative first Chern class. Letting $`\delta `$ be the class of the diagonal we see that $`f_i^2=0`$, $`f_1.f_2=1`$, $`f_i.\delta =1`$ and $`\delta ^2=22g`$.
Let $`f=f_1+f_2`$ and for convenience make the change of variables $`\delta ^{}=\delta f`$. Then we have the following intersection numbers on $`X`$:
$$f^2=2,\delta ^{}.f=0,\text{and}\delta ^2=2g.$$
Now consider the $``$-divisor
$$L_t=tf\delta ^{}$$
which is ample for $`t0`$. We define
$$s_C=\text{inf}\{t:L_t\text{ is ample}\}.$$
Clearly $`s_C\sqrt{g}`$ for if $`L_t`$ is ample then $`0<L_t^2=2t^22g`$. In fact conjecturally $`s_C=\sqrt{g}`$ for “most” curves (see Remark 4.5).
Recall that a branched cover from a curve to $`^1`$ is said to be simple if it has only ramifications which are locally $`zz^2`$ and no two ramification points map to the same point in $`^1`$. Given a simple branched cover $`\pi :C^1`$ of degree $`d`$, consider the fibre product $`C\times _\pi CX`$. This contains the diagonal $`\mathrm{\Delta }`$ and we let $`Z`$ be the residual divisor so $`Z+\mathrm{\Delta }=C\times _\pi C`$.
###### Lemma 3.1.
The class of $`Z`$ in $`N^1(X)_{}`$ is $`(d1)f\delta ^{}`$, so $`Z`$ has self-intersection $`Z^2=2(d1)^22g.`$
###### Proof.
In the product $`^1\times ^1`$ let $`F_i`$, $`i=1,2`$ be the numerical class of the coordinate planes and $`D`$ be the class of the diagonal. Then $`D=F_1+F_2`$ so
$$Z+\delta =C\times _\pi C=(\pi \times \pi )^{}(D)=(\pi \times \pi )^{}(F_1+F_2)=df_1+df_2=df.$$
Hence $`Z=df\delta =(d1)f\delta ^{}`$ as claimed. ∎
We will be interested in the case when $`2d1<\sqrt{g}`$. Then $`Z^2`$ is negative, and we will show that $`X`$ is not slope semistable with respect to $`Z`$ for suitable polarisations.
###### Theorem 3.2 (Kouvidakis ).
Suppose $`C`$ is a smooth curve of genus $`g2`$ which admits a simple branched cover to $`^1`$ of degree $`d`$ with $`d1\sqrt{g}`$. Then $`s_C=\frac{g}{d1}`$.
###### Proof.
See Theorem 1.5.8.
###### Theorem 3.3.
Suppose $`C`$ is a smooth curve of genus $`g`$ which admits a simple branched cover $`\pi :C^1`$ of degree $`d`$ with $`2d1<\sqrt{g}`$.
Then $`X=C\times C`$ is not slope semistable with respect $`L_t`$ for $`t`$ sufficiently close to $`s_C`$.
###### Proof.
Let $`t>s_C`$ so $`L_t`$ is ample. The canonical divisor of $`X`$ is $`K=(2g2)f`$ so from (2.4)
$$\mu (X,L_t)=\frac{K.L_t}{L_t^2}=\frac{t(2g2)}{t^2g}.$$
(3.4)
By (3.2), $`s_C=\frac{g}{d1}`$. Letting $`Z`$ be the curve from Lemma 3.1 whose class is $`(d1)f\delta ^{}`$ we now bound the Seshadri constant of $`Z`$. Since
$$t(d1)>s_C(d1)=\frac{g}{d1}(d1)>0$$
we have that $`L_tZ=tf\delta ^{}((d1)f\delta ^{})=(t(d1))f`$ is ample. Thus $`ϵ(Z,L_t)1`$.
To calculate the slope of $`Z`$ we need the quantities
$`L_t.Z`$ $`=`$ $`(tf\delta ^{}).((d1)f\delta ^{})`$
$`=`$ $`2t(d1)2g,`$
$`K.Z`$ $`=`$ $`(2g2)f.((d1)f\delta ^{})=2(2g2)(d1),`$
$`Z^2`$ $`=`$ $`((d1)f\delta ^{})^2=2(d1)^22g.`$
Thus from (2.4),
$`\mu _1(𝒪_Z,L_t)`$ $`=`$ $`{\displaystyle \frac{3(2L_t.Z(K.Z+Z^2))}{2(3L_t.ZZ^2)}}`$
$`=`$ $`{\displaystyle \frac{3(4t(d1)4g2(2g2)(d1)2(d1)^2+2g)}{2(6t(d1)6g2(d1)^2+2g)}}.`$
We claim that $`\mu _1(𝒪_Z,L_t)<\mu (X,L_t)`$ as $`t`$ tends to $`s_C=\frac{g}{d1}`$ from above. Since this is an open condition it is sufficient to show that it holds when $`t=s_C`$. By (3.4, 3),
$`\mu (X,L_{s_C})`$ $`=`$ $`{\displaystyle \frac{(2g2)(d1)}{g(d1)^2}},`$
$`\mu _1(𝒪_Z,L_{s_C})`$ $`=`$ $`{\displaystyle \frac{3(g(2g2)(d1)(d1)^2)}{2(g(d1)^2)}}`$
(notice that our assumption $`d1<\sqrt{g}`$ ensures that both of these are finite). Hence as $`d12`$,
$`2(g(d1)^2)`$ $`.(\mu _1(𝒪_Z,L_{s_C})\mu (X,L_{s_C}))`$
$`=3g3(2g2)(d1)3(d1)^2+2(2g2)(d1)`$
$`=3g(2g2)(d1)3(d1)^2`$
$`3g2(2g2)12<0.`$
Thus $`ϵ(Z,L_t)1`$ and $`\mu _1(𝒪_Z,L_t)<\mu (X,L_t)`$ as $`t`$ tends to $`s_C`$ from above, which proves that $`(X,L_t)`$ is not slope semistable. ∎
###### Theorem 3.7.
For $`g5`$ there exist smooth curves $`C`$ of genus $`g`$ such that $`X=C\times C`$ is not slope semistable with respect to certain polarisations. Thus there are Kähler classes on $`X`$ that do not admit Kähler metrics of constant scalar curvature.
###### Proof.
By the Riemann existence theorem there exist smooth curves of genus $`g`$ which admit a simple branched covering over $`^1`$ of degree $`2d1<\sqrt{g}`$ (in fact one can even take $`d=3`$). So by (3.3), $`X=C\times C`$ is not slope stable with respect to certain polarisations and thus not K-semistable by (2.3). The application to cscK metrics comes from (2.1). ∎
###### Theorem 3.8.
For $`g5`$ there exist smooth curves $`C`$ of genus $`g`$ such that $`X=C\times C`$ is not asymptotically Hilbert semistable (resp. not asymptotically Chow semistable) with respect to suitable polarisations.
###### Proof.
This follows from (3.7) and (2.6). ∎
###### Remark 3.9.
Let $`(X_i,L_i)`$ $`i=1,2`$ be polarised manifolds and $`\pi _i:X_1\times X_2X_i`$ be the projection maps. If $`L=\pi _1^{}L_1+\pi _2^{}L_2`$ then
$$\mu (X_1\times X_2,L)=\mu (X_1,L_1)+\mu (X_2,L_2).$$
Moreover if $`Z`$ is subscheme of $`X_1`$ then one can calculate the slope of $`Z\times X_2X_1\times X_2`$ is $`\mu _c(𝒪_{Z\times X_1},L)=\mu _c(𝒪_Z,L_1)+\mu (X_2,L_2)`$. Thus if $`(X_1,L_1)`$ is slope unstable so is the product $`(X_1\times X_2,L)`$.
So by taking the product of an slope unstable surface with any manifold with negative first Chern class, we get manifolds of any dimension $`n2`$ with negative first Chern class which have Kähler classes that do not admit cscK metrics.
## 4. The J-flow on products of curves
The Mabuchi functional for a given Kähler class $`\mathrm{\Omega }`$ on a complex manifold $`X`$ has as its critical points the metrics which are cscK. Conjecturally the existence of a cscK metric in $`\mathrm{\Omega }`$ is equivalent to the properness of the Mabuchi functional. This is known to be true when $`\mathrm{\Omega }`$ is proportional to the canonical class ; and when $`\mathrm{\Omega }`$ admits a cscK metric the Mabuchi functional is necessarily bounded from below
Now suppose that $`X`$ has negative first Chern class. Chen introduces a flow on Kähler manifolds, called the J-flow, and points out that convergence of this flow implies lower boundedness of the Mabuchi functional. In it is shown that on a surface with negative first Chern class, the J-flow converges as long as the class $`2(_Xc_1(X).\mathrm{\Omega })\mathrm{\Omega }+\left(_X\mathrm{\Omega }^2\right)c_1(X)`$ is positive.
The following theorem is a particular case of Theorems 1.2 and 1.4 in applied to polarised surfaces.
###### Theorem 4.1.
Let $`(X,L)`$ be a polarised surface and define a divisor by
$$\alpha =2(K.L)L(L^2)K.$$
* If $`\alpha `$ is ample then the J-flow converges and the Mabuchi functional is proper on the class $`c_1(L)`$.
* If $`\alpha `$ is not ample then there exist $`m`$ irreducible curves $`E_i`$ of negative self intersection and positive numbers $`a_i`$ such that
$$\alpha \underset{i=1}{\overset{m}{}}a_iE_i\text{ is ample}.$$
(4.2)
In fact if $`\alpha `$ is not ample then the J-flow “blows up” along the intersection of all divisors in the linear series $`|\mathrm{\Sigma }a_iE_i|`$ .
From this theorem one might expect that $`(X,L)`$ is stable if $`\alpha `$ is ample, and perhaps that if $`\alpha `$ is not ample that the $`E_i`$ witness instability of $`(X,L)`$ (c.f. Remark 4.7).
When $`C`$ is a curve, $`X=C\times C`$ and $`L=L_t=tf\delta ^{}`$ with $`t>s_C`$ it is easy to determine when $`\alpha `$ is ample.
###### Lemma 4.3.
Let $`L=L_t=tf\delta ^{}`$ with $`t>s_C`$. Then $`\alpha =2(K.L)L(L^2)K`$ is ample if and only if $`t^2+g>2ts_C`$ if and only if $`t>s_C+\sqrt{s_C^2g}`$.
###### Proof.
As $`L_t^2=2t^22g`$ and $`K.L_t=2t(2g2)`$ we have
$`\alpha `$ $`=`$ $`4t(2g2)(tf\delta ^{})(2t^22g)(2g2)f`$
$`=`$ $`2(2g2)((t^2+g)f2t\delta ^{})`$
which is ample if and only if $`t^2+g>2ts_C`$ which occurs if and only if $`h(t)=t^22ts_C+g>0`$. But the roots of $`h`$ are $`s_C\pm \sqrt{s_C^2g}`$ and $`s_C\sqrt{g}`$ so the lemma follows. ∎
Thus one can deduce properness of the Mabuchi functional when $`s_C=\sqrt{g}`$.
###### Corollary 4.4.
If $`s_C=\sqrt{g}`$ then the Mabuchi function is proper on any class on $`X=C\times C`$ of the form $`c_1(L_t)`$ for $`t>s_C`$.
###### Remark 4.5.
If $`g`$ is a perfect square and $`C`$ is very general curve of genus $`g`$ then $`s_C=\sqrt{g}`$ ( Corollary 1.5.9). If the Nagata conjecture holds and $`g10`$ then the same conclusion holds without the hypothesis that $`g`$ is a perfect square .
###### Remark 4.6.
Let $`C`$ be as in (3.3) and $`X=C\times C`$. It would be interesting to know if the Mabuchi functional is bounded on the class $`c_1(L_t)`$ if and only if $`t>s_C+\sqrt{s_C^2g}`$.
It is not the case that the curve $`Z=(d1)f\delta ^{}`$ from (3.3) always slope destabilises when $`t<s_C+\sqrt{s_C^2g}`$. For example let $`C`$ be a curve of genus $`g=5`$ admitting a simple branched cover to $`^1`$ of degree $`d=3`$. Then $`s_C=5/2`$ so $`s_C+\sqrt{s_C^2g}=\frac{5+\sqrt{5}}{2}>3`$. Put $`t=3`$ so $`L_3=3f\delta ^{}`$ is ample. The slope of $`Z`$ is
$$\mu _c(𝒪_Z,L)=\frac{3(215c)}{2c(3+c)}$$
which is greater than $`\mu (X,L)=6`$ for all $`c>0`$.
By the previous lemma $`\alpha `$ is not ample. It is unfortunately not the case that any divisor satisfying (4.2) necessarily destabilise (c.f. Remark 4.7) because in the example above, $`\alpha 33Z=158f63\delta ^{}`$ is ample.
jaross@math.columbia.edu
Department of Mathematics, Columbia University, New York, NY 10027. USA. |
warning/0506/quant-ph0506216.html | ar5iv | text | # Probabilistic teleportation of unknown two-particle state via POVM
## Abstract
We propose a scheme for probabilistic teleportation of unknown two-particle state with partly entangled four-particle state via POVM. In this scheme the teleportation of unknown two-particle state can be realized with certain probability by performing two Bell state measurements, a proper POVM and a unitary transformation.
Quantum teleportation is a process of transmission of an unknown quantum state via a previously shared EPR pair with the help of only two classical bits transmitted through a classical channel Bennett93 . It was regarded as one of the most striking progress of quantum information theory NielsenChuang . It may have a number of useful applications in quantum computer CiracZoller ; Barenco , quantum dense coding BennenttWiesner , quantum cryptography BB84 ; Ekert ; Bennett92 and quantum secure direct communication ShimizuImoto ; Beige ; DengLong ; YanZhang ; Gao .
Since Bennett et al. showed that an unknown quantum state of two-state particle (or qubit) can be teleported from a sender Alice to a spatially distant receiver Bob in 1993 Bennett93 , people have paid much attention to quantum teleportation. Research work on quantum teleportation was soon widely started up, and has got great development, theoretical and experimental as well. On the one hand, the teleportation of a photon polarization state has been demonstrated experimentally with the use of polarization-entangled photons BPZ and path-entangled photons Boschi , respectively. The teleportation of a coherent state corresponding to continuous variable system was also realized in the laboratory Furusawa . On the other hand, the idea of quantum teleportation has been generalized to many cases GorbachevTrubilko ; LuGuo ; Bandyopadhyay ; Rigolin ; ShiGuo ; ZengLong ; MorHorodecki ; LiuGuo ; DaiLi ; YanTanYang ; YanWang ; YanYang ; YanBai ; GaoYanWang ; GaoWangYanChinPhysLett ; GaoYanWangChinesePhysics ; Gao2 .
Mor and Horodecki discussed the problems of teleportation via positive operator valued measure (POVM) Helstrom , which was also called generalized measurement, and conclusive teleportation MorHorodecki . It was showed that a perfect conclusive teleportation can be obtained with any pure entangled state. Bandyopadhyay presented two optimal methods of teleporting an unknown qubit using any pure entangled state Bandyopadhyay . About three years ago we designed a scheme for probabilistic teleporting an unknown two-particle state with partly pure entangled four-particle state via a projective measurement on an auxiliary particle YanTanYang . Is it possible to use POVM in the probabilistic teleportation of an unknown two-particle state with partly pure entangled four-particle state? In this Letter, we will try to answer this question, i.e., to propose a scheme for probabilistic teleportation of unknown two-particle state with partly entangled four-particle state via POVM. It will be shown that by performing two Bell state measurements, a proper POVM and a unitary transformation, the unknown two-particle state can be teleported from the sender Alice to the receiver Bob with certain probability. We state our scheme in details as follows.
Suppose that the sender Alice has two particles 1,2 in an unknown state
$$|\mathrm{\Phi }_{12}=(a|00+b|01+c|10+d|11)_{12},$$
(1)
where $`a,b,c,d`$ are arbitrary complex numbers, and satisfy $`|a|^2+|b|^2+|c|^2+|d|^2=1`$. We also suppose that Alice and Bob share quantum entanglement in the form of following partly pure entangled four-particle state, which will be used as the quantum channel,
$$|\mathrm{\Phi }_{3456}=(\alpha |0000+\beta |1001+\gamma |0110+\delta |1111)_{3456},$$
(2)
where $`\alpha ,\beta ,\gamma ,\delta `$ are nonzero real numbers, and $`\alpha ^2+\beta ^2+\gamma ^2+\delta ^2=1`$. The particles 3 and 4, and particle pair (1,2) are in Alice’s possession, and other two particles 5 and 6 are in Bob’s possession. The overall state of six particles is
$$|\mathrm{\Phi }_w=|\mathrm{\Phi }_{12}|\mathrm{\Phi }_{3456}.$$
(3)
In order to realize the teleportation, firstly Alice performs two Bell state measurements on particles 2,3 and 1,4, then the resulting state of Bob’s particles 5,6 will be one of the following states YanTanYang
$$|\mathrm{\Psi }_0_{56}=\frac{{}_{14}{}^{}\mathrm{\Phi }^+|_{23}\mathrm{\Phi }^+|\mathrm{\Phi }_w}{|_{14}\mathrm{\Phi }^+|_{23}\mathrm{\Phi }^+|\mathrm{\Phi }_w|}=\frac{1}{\sqrt{|a\alpha |^2+|b\beta |^2+|c\gamma |^2+|d\delta |^2}}(a\alpha |00+b\beta |01+c\gamma |10+d\delta |11)_{56},$$
(4)
$$|\mathrm{\Psi }_1_{56}=\frac{{}_{14}{}^{}\mathrm{\Phi }^{}|_{23}\mathrm{\Phi }^+|\mathrm{\Phi }_w}{|_{14}\mathrm{\Phi }^{}|_{23}\mathrm{\Phi }^+|\mathrm{\Phi }_w|}=\frac{1}{\sqrt{|a\alpha |^2+|b\beta |^2+|c\gamma |^2+|d\delta |^2}}(a\alpha |00+b\beta |01c\gamma |10d\delta |11)_{56},$$
(5)
$$|\mathrm{\Psi }_2_{56}=\frac{{}_{14}{}^{}\mathrm{\Psi }^+|_{23}\mathrm{\Phi }^+|\mathrm{\Phi }_w}{|_{14}\mathrm{\Psi }^+|_{23}\mathrm{\Phi }^+|\mathrm{\Phi }_w|}=\frac{1}{\sqrt{|a\gamma |^2+|b\delta |^2+|c\alpha |^2+|d\beta |^2}}(a\gamma |10+b\delta |11+c\alpha |00+d\beta |01)_{56},$$
(6)
$$|\mathrm{\Psi }_3_{56}=\frac{{}_{14}{}^{}\mathrm{\Psi }^+|_{23}\mathrm{\Phi }^{}|\mathrm{\Phi }_w}{|_{14}\mathrm{\Psi }^+|_{23}\mathrm{\Phi }^{}|\mathrm{\Phi }_w|}=\frac{1}{\sqrt{|a\gamma |^2+|b\delta |^2+|c\alpha |^2+|d\beta |^2}}(a\gamma |10+b\delta |11c\alpha |00d\beta |01)_{56},$$
(7)
$$|\mathrm{\Psi }_4_{56}=\frac{{}_{14}{}^{}\mathrm{\Phi }^+|_{23}\mathrm{\Phi }^{}|\mathrm{\Phi }_w}{|_{14}\mathrm{\Phi }^+|_{23}\mathrm{\Phi }^{}|\mathrm{\Phi }_w|}=\frac{1}{\sqrt{|a\alpha |^2+|b\beta |^2+|c\gamma |^2+|d\delta |^2}}(a\alpha |00b\beta |01+c\gamma |10d\delta |11)_{56},$$
(8)
$$|\mathrm{\Psi }_5_{56}=\frac{{}_{14}{}^{}\mathrm{\Phi }^{}|_{23}\mathrm{\Phi }^{}|\mathrm{\Phi }_w}{|_{14}\mathrm{\Phi }^{}|_{23}\mathrm{\Phi }^{}|\mathrm{\Phi }_w|}=\frac{1}{\sqrt{|a\alpha |^2+|b\beta |^2+|c\gamma |^2+|d\delta |^2}}(a\alpha |00b\beta |01c\gamma |10+d\delta |11)_{56},$$
(9)
$$|\mathrm{\Psi }_6_{56}=\frac{{}_{14}{}^{}\mathrm{\Psi }^+|_{23}\mathrm{\Phi }^{}|\mathrm{\Phi }_w}{|_{14}\mathrm{\Psi }^+|_{23}\mathrm{\Phi }^{}|\mathrm{\Phi }_w|}=\frac{1}{\sqrt{|a\gamma |^2+|b\delta |^2+|c\alpha |^2+|d\beta |^2}}(a\gamma |10b\delta |11+c\alpha |00d\beta |01)_{56},$$
(10)
$$|\mathrm{\Psi }_7_{56}=\frac{{}_{14}{}^{}\mathrm{\Psi }^{}|_{23}\mathrm{\Phi }^{}|\mathrm{\Phi }_w}{|_{14}\mathrm{\Psi }^{}|_{23}\mathrm{\Phi }^{}|\mathrm{\Phi }_w|}=\frac{1}{\sqrt{|a\gamma |^2+|b\delta |^2+|c\alpha |^2+|d\beta |^2}}(a\gamma |10b\delta |11c\alpha |00+d\beta |01)_{56},$$
(11)
$$|\mathrm{\Psi }_8_{56}=\frac{{}_{14}{}^{}\mathrm{\Phi }^+|_{23}\mathrm{\Psi }^+|\mathrm{\Phi }_w}{|_{14}\mathrm{\Phi }^+|_{23}\mathrm{\Psi }^+|\mathrm{\Phi }_w|}=\frac{1}{\sqrt{|a\beta |^2+|b\alpha |^2+|c\delta |^2+|d\gamma |^2}}(a\beta |01+b\alpha |00+c\delta |11+d\gamma |10)_{56},$$
(12)
$$|\mathrm{\Psi }_9_{56}=\frac{{}_{14}{}^{}\mathrm{\Phi }^{}|_{23}\mathrm{\Psi }^+|\mathrm{\Phi }_w}{|_{14}\mathrm{\Phi }^{}|_{23}\mathrm{\Psi }^+|\mathrm{\Phi }_w|}=\frac{1}{\sqrt{|a\beta |^2+|b\alpha |^2+|c\delta |^2+|d\gamma |^2}}(a\beta |01+b\alpha |00c\delta |11d\gamma |10)_{56},$$
(13)
$$|\mathrm{\Psi }_{10}_{56}=\frac{{}_{14}{}^{}\mathrm{\Psi }^+|_{23}\mathrm{\Psi }^+|\mathrm{\Phi }_w}{|_{14}\mathrm{\Psi }^+|_{23}\mathrm{\Psi }^+|\mathrm{\Phi }_w|}=\frac{1}{\sqrt{|a\delta |^2+|b\gamma |^2+|c\beta |^2+|d\alpha |^2}}(a\delta |11+b\gamma |10+c\beta |01+d\alpha |00)_{56},$$
(14)
$$|\mathrm{\Psi }_{11}_{56}=\frac{{}_{14}{}^{}\mathrm{\Psi }^{}|_{23}\mathrm{\Psi }^+|\mathrm{\Phi }_w}{|_{14}\mathrm{\Psi }^{}|_{23}\mathrm{\Psi }^+|\mathrm{\Phi }_w|}=\frac{1}{\sqrt{|a\delta |^2+|b\gamma |^2+|c\beta |^2+|d\alpha |^2}}(a\delta |11+b\gamma |10c\beta |01d\alpha |00)_{56},$$
(15)
$$|\mathrm{\Psi }_{12}_{56}=\frac{{}_{14}{}^{}\mathrm{\Phi }^+|_{23}\mathrm{\Psi }^{}|\mathrm{\Phi }_w}{|_{14}\mathrm{\Phi }^+|_{23}\mathrm{\Psi }^{}|\mathrm{\Phi }_w|}=\frac{1}{\sqrt{|a\beta |^2+|b\alpha |^2+|c\delta |^2+|d\gamma |^2}}(a\beta |01b\alpha |00+c\delta |11d\gamma |10)_{56},$$
(16)
$$|\mathrm{\Psi }_{13}_{56}=\frac{{}_{14}{}^{}\mathrm{\Phi }^{}|_{23}\mathrm{\Psi }^{}|\mathrm{\Phi }_w}{|_{14}\mathrm{\Phi }^{}|_{23}\mathrm{\Psi }^{}|\mathrm{\Phi }_w|}=\frac{1}{\sqrt{|a\beta |^2+|b\alpha |^2+|c\delta |^2+|d\gamma |^2}}(a\beta |01b\alpha |00c\delta |11+d\gamma |10)_{56},$$
(17)
$$|\mathrm{\Psi }_{14}_{56}=\frac{{}_{14}{}^{}\mathrm{\Psi }^+|_{23}\mathrm{\Psi }^{}|\mathrm{\Phi }_w}{|_{14}\mathrm{\Psi }^+|_{23}\mathrm{\Psi }^{}|\mathrm{\Phi }_w|}=\frac{1}{\sqrt{|a\delta |^2+|b\gamma |^2+|c\beta |^2+|d\alpha |^2}}(a\delta |11b\gamma |10+c\beta |01d\alpha |00)_{56},$$
(18)
$$|\mathrm{\Psi }_{15}_{56}=\frac{{}_{14}{}^{}\mathrm{\Psi }^{}|_{23}\mathrm{\Psi }^{}|\mathrm{\Phi }_w}{|_{14}\mathrm{\Psi }^{}|_{23}\mathrm{\Psi }^{}|\mathrm{\Phi }_w|}=\frac{1}{\sqrt{|a\delta |^2+|b\gamma |^2+|c\beta |^2+|d\alpha |^2}}(a\delta |11b\gamma |10c\beta |01+d\alpha |00)_{56}.$$
(19)
Here
$$\begin{array}{ccc}|\mathrm{\Phi }^\pm =\frac{1}{\sqrt{2}}(|00\pm |11),& & |\mathrm{\Psi }^\pm =\frac{1}{\sqrt{2}}(|01\pm |10)\end{array}$$
(20)
are the Bell states. Then Alice informs Bob her two Bell state measurement outcomes via a classical channel. By outcomes received, Bob can determine the state of particles 5,6 exactly. Without loss of generality, next we will give the case for $`|\mathrm{\Psi }_0_{56}`$, the other cases can be deduced similarly. In order to realize the teleportation, Bob introduces two auxiliary qubits $`a,b`$ in the state $`|00_{ab}`$. So the state of particles $`5,6,a,b`$ becomes
$$|\mathrm{\Psi }_0_{56}|00_{ab}=\frac{1}{\sqrt{|a\alpha |^2+|b\beta |^2+|c\gamma |^2+|d\delta |^2}}(a\alpha |00+b\beta |01+c\gamma |10+d\delta |11)_{56}|00_{ab}.$$
(21)
Then Bob performs two controlled-not operations with particles 5,6 as the controlled qubits and the auxiliary particles $`a,b`$ as the target qubits respectively. After completing this operation the particles $`5,6,a,b`$ are in state
$$|\mathrm{\Psi }_0^{}_{56ab}=\frac{1}{\sqrt{|a\alpha |^2+|b\beta |^2+|c\gamma |^2+|d\delta |^2}}(a\alpha |0000+b\beta |0101+c\gamma |1010+d\delta |1111)_{56ab}.$$
(22)
A simple algebraic rearrangement of this expression yields
$$\begin{array}{cc}|\mathrm{\Psi }_0^{}_{56ab}=& \frac{1}{4\sqrt{|a\alpha |^2+|b\beta |^2+|c\gamma |^2+|d\delta |^2}}[(a|00+b|01+c|10+d|11)_{56}(\alpha |00+\beta |01+\gamma |10+\delta |11)_{ab}\hfill \\ & +(a|00+b|01c|10d|11)_{56}(\alpha |00+\beta |01\gamma |10\delta |11)_{ab}\hfill \\ & +(a|00b|01+c|10d|11)_{56}(\alpha |00\beta |01+\gamma |10\delta |11)_{ab}\hfill \\ & +(a|00b|01c|10+d|11)_{56}(\alpha |00\beta |01\gamma |10+\delta |11)_{ab}].\hfill \end{array}$$
(23)
At this stage, Bob makes an optimal POVM Bandyopadhyay ; Helstrom on the ancillary particle $`a,b`$ to conclusively distinguish the above states. We choose the optimal POVM in this subspace as follows
$$\begin{array}{ccccc}P_1=\frac{1}{x}|\mathrm{\Psi }_1\mathrm{\Psi }_1|,& P_2=\frac{1}{x}|\mathrm{\Psi }_2\mathrm{\Psi }_2|,& P_3=\frac{1}{x}|\mathrm{\Psi }_3\mathrm{\Psi }_3|,& P_4=\frac{1}{x}|\mathrm{\Psi }_4\mathrm{\Psi }_4|,& P_5=I\frac{1}{x}\mathrm{\Sigma }_{i=1}^4|\mathrm{\Psi }_i\mathrm{\Psi }_i|,\end{array}$$
(24)
where
$$|\mathrm{\Psi }_1=\frac{1}{\sqrt{\frac{1}{\alpha ^2}+\frac{1}{\beta ^2}+\frac{1}{\gamma ^2}+\frac{1}{\delta ^2}}}(\frac{1}{\alpha }|00+\frac{1}{\beta }|01+\frac{1}{\gamma }|10+\frac{1}{\delta }|11)_{ab},$$
(25)
$$|\mathrm{\Psi }_2=\frac{1}{\sqrt{\frac{1}{\alpha ^2}+\frac{1}{\beta ^2}+\frac{1}{\gamma ^2}+\frac{1}{\delta ^2}}}(\frac{1}{\alpha }|00+\frac{1}{\beta }|01\frac{1}{\gamma }|10\frac{1}{\delta }|11)_{ab},$$
(26)
$$|\mathrm{\Psi }_3=\frac{1}{\sqrt{\frac{1}{\alpha ^2}+\frac{1}{\beta ^2}+\frac{1}{\gamma ^2}+\frac{1}{\delta ^2}}}(\frac{1}{\alpha }|00\frac{1}{\beta }|01+\frac{1}{\gamma }|10\frac{1}{\delta }|11)_{ab},$$
(27)
$$|\mathrm{\Psi }_4=\frac{1}{\sqrt{\frac{1}{\alpha ^2}+\frac{1}{\beta ^2}+\frac{1}{\gamma ^2}+\frac{1}{\delta ^2}}}(\frac{1}{\alpha }|00\frac{1}{\beta }|01\frac{1}{\gamma }|10+\frac{1}{\delta }|11)_{ab};$$
(28)
$`I`$ is an identity operator; $`x`$ is a coefficient relating to $`\alpha ,\beta ,\gamma ,\delta `$, $`1x4`$, and makes $`P_5`$ to be a positive operator. For exactly determining $`x`$, we would like to write the five operators $`P_1,P_2,P_3,P_4,P_5`$ in the matrix form
$$P_1=\frac{1}{x(\frac{1}{\alpha ^2}+\frac{1}{\beta ^2}+\frac{1}{\gamma ^2}+\frac{1}{\delta ^2})}\left(\begin{array}{cccc}\frac{1}{\alpha ^2}& \frac{1}{\alpha \beta }& \frac{1}{\alpha \gamma }& \frac{1}{\alpha \delta }\\ \frac{1}{\alpha \beta }& \frac{1}{\beta ^2}& \frac{1}{\beta \gamma }& \frac{1}{\beta \delta }\\ \frac{1}{\alpha \gamma }& \frac{1}{\beta \gamma }& \frac{1}{\gamma ^2}& \frac{1}{\gamma \delta }\\ \frac{1}{\alpha \delta }& \frac{1}{\beta \delta }& \frac{1}{\gamma \delta }& \frac{1}{\delta ^2}\end{array}\right),$$
(29)
$$P_2=\frac{1}{x(\frac{1}{\alpha ^2}+\frac{1}{\beta ^2}+\frac{1}{\gamma ^2}+\frac{1}{\delta ^2})}\left(\begin{array}{cccc}\frac{1}{\alpha ^2}& \frac{1}{\alpha \beta }& \frac{1}{\alpha \gamma }& \frac{1}{\alpha \delta }\\ \frac{1}{\alpha \beta }& \frac{1}{\beta ^2}& \frac{1}{\beta \gamma }& \frac{1}{\beta \delta }\\ \frac{1}{\alpha \gamma }& \frac{1}{\beta \gamma }& \frac{1}{\gamma ^2}& \frac{1}{\gamma \delta }\\ \frac{1}{\alpha \delta }& \frac{1}{\beta \delta }& \frac{1}{\gamma \delta }& \frac{1}{\delta ^2}\end{array}\right),$$
(30)
$$P_3=\frac{1}{x(\frac{1}{\alpha ^2}+\frac{1}{\beta ^2}+\frac{1}{\gamma ^2}+\frac{1}{\delta ^2})}\left(\begin{array}{cccc}\frac{1}{\alpha ^2}& \frac{1}{\alpha \beta }& \frac{1}{\alpha \gamma }& \frac{1}{\alpha \delta }\\ \frac{1}{\alpha \beta }& \frac{1}{\beta ^2}& \frac{1}{\beta \gamma }& \frac{1}{\beta \delta }\\ \frac{1}{\alpha \gamma }& \frac{1}{\beta \gamma }& \frac{1}{\gamma ^2}& \frac{1}{\gamma \delta }\\ \frac{1}{\alpha \delta }& \frac{1}{\beta \delta }& \frac{1}{\gamma \delta }& \frac{1}{\delta ^2}\end{array}\right),$$
(31)
$$P_4=\frac{1}{x(\frac{1}{\alpha ^2}+\frac{1}{\beta ^2}+\frac{1}{\gamma ^2}+\frac{1}{\delta ^2})}\left(\begin{array}{cccc}\frac{1}{\alpha ^2}& \frac{1}{\alpha \beta }& \frac{1}{\alpha \gamma }& \frac{1}{\alpha \delta }\\ \frac{1}{\alpha \beta }& \frac{1}{\beta ^2}& \frac{1}{\beta \gamma }& \frac{1}{\beta \delta }\\ \frac{1}{\alpha \gamma }& \frac{1}{\beta \gamma }& \frac{1}{\gamma ^2}& \frac{1}{\gamma \delta }\\ \frac{1}{\alpha \delta }& \frac{1}{\beta \delta }& \frac{1}{\gamma \delta }& \frac{1}{\delta ^2}\end{array}\right),$$
(32)
$$\begin{array}{c}P_5=\frac{1}{\frac{1}{\alpha ^2}+\frac{1}{\beta ^2}+\frac{1}{\gamma ^2}+\frac{1}{\delta ^2}}\hfill \\ \left(\begin{array}{cccc}(1\frac{4}{x})\frac{1}{\alpha ^2}+\frac{1}{\beta ^2}+\frac{1}{\gamma ^2}+\frac{1}{\delta ^2}& 0& 0& 0\\ 0& (1\frac{4}{x})\frac{1}{\beta ^2}+\frac{1}{\alpha ^2}+\frac{1}{\gamma ^2}+\frac{1}{\delta ^2}& 0& 0\\ 0& 0& (1\frac{4}{x})\frac{1}{\gamma ^2}+\frac{1}{\alpha ^2}+\frac{1}{\beta ^2}+\frac{1}{\delta ^2}& 0\\ 0& 0& 0& (1\frac{4}{x})\frac{1}{\delta ^2}+\frac{1}{\alpha ^2}+\frac{1}{\beta ^2}+\frac{1}{\gamma ^2}\end{array}\right).\hfill \end{array}$$
(33)
Obviously, we should carefully choose $`x`$ such that all the diagonal elements of $`P_5`$ are nonnegative. If the result of Bob’s POVM is $`P_1`$, then Bob can safely conclude that the state of the particles 5,6 is
$$|\mathrm{\Phi }_{56}=(a|00+b|01+c|10+d|11)_{56}.$$
(34)
If Bob’s POVM outcome is $`P_2`$, Bob can obtain $`|\mathrm{\Phi }_{56}`$ by applying the unitary transformation $`\sigma _zI`$ on the particles 5,6. If the Bob’s POVM outcome $`P_3`$ occurs, Bob makes $`I\sigma _z`$ on the particles 5,6 to recover $`|\mathrm{\Phi }_{56}`$. If Bob’s POVM outcome is $`P_4`$, Bob applies $`\sigma _z\sigma _z`$ on the particles 5,6 to recover $`|\mathrm{\Phi }_{56}`$. Therefore, in these four cases stated above, the teleportation is realized successfully. However if Bob’s POVM outcome is $`P_5`$, Bob can infer nothing about the identity of the state of the particles 5,6. In this case the teleportation fails. As a matter of fact, the key point is that Bob never makes a mistake identifying the state of the particles 5,6. This infallibility comes at the price that sometime Bob obtains no information about the identity of the state of the particles 5,6.
By the similar method we can make the teleportation successful in the other outcomes of Alice’s Bell state measurement. For the sake of saving the space we will not write them out.
Evidently, when the Bell states $`|\mathrm{\Phi }^+_{23}`$ and $`|\mathrm{\Phi }^+_{14}`$ are acquired in Alice’s two Bell state measurements, the probability of successful teleportation is
$$\begin{array}{cc}& |_{14}\mathrm{\Phi }^+|_{23}\mathrm{\Phi }^+|\mathrm{\Phi }_w|^2(1_{56ab}\mathrm{\Psi }_0^{}|P_5I|\mathrm{\Psi }_0^{}_{56ab})\hfill \\ \hfill =& \frac{|a\alpha |^2+|b\beta |^2+|c\gamma |^2+|d\delta |^2}{4}\times \frac{4}{x(\frac{1}{\alpha ^2}+\frac{1}{\beta ^2}+\frac{1}{\gamma ^2}+\frac{1}{\delta ^2})(|a\alpha |^2+|b\beta |^2+|c\gamma |^2+|d\delta |^2)}\hfill \\ \hfill =& \frac{1}{x(\frac{1}{\alpha ^2}+\frac{1}{\beta ^2}+\frac{1}{\gamma ^2}+\frac{1}{\delta ^2})}.\hfill \end{array}$$
(35)
Synthesizing all Alice’s Bell state measurement cases (sixteen kinds in all), the probability of successful teleportation in this scheme is
$$p=\frac{16}{x(\frac{1}{\alpha ^2}+\frac{1}{\beta ^2}+\frac{1}{\gamma ^2}+\frac{1}{\delta ^2})}.$$
(36)
Apparently, if the quantum channel is made up of the maximum entangled state $`\frac{1}{2}(|0000+|1001+|0110+|1111)_{3456}`$, i.e. $`\alpha =\beta =\gamma =\delta =\frac{1}{2}`$, we can choose $`x=1`$, then $`P_5`$ is zero operator. In this case the probabilistic teleportation becomes usual teleportation.
In summary, a scheme for probabilistic teleporting the unknown quantum state of two-particle is proposed. We hope that this scheme will be realized by experiment.
###### Acknowledgements.
This work was supported by Hebei Natural Science Foundation of China under Grant No: A2004000141 and A2005000140, and Natural Science Foundation of Hebei Normal University. |
warning/0506/astro-ph0506066.html | ar5iv | text | # The Star Cluster Population of M51: III. Cluster disruption and formation history
## 1 Introduction
The goal of this series of papers is to understand the properties of the entire star cluster population of the interacting spiral galaxy M51. These properties include the age and mass distribution of the cluster population. Additional properties are the survival rate of the clusters as well as any relations between the observed properties. These relations may be used to constrain cluster formation and destruction scenarios.
In order to study the above properties, we exploit the large amount of HST broad-band archival data on M51, which covers roughly 50% of the observed surface area of M51, and covers a broad spectral range (UV to NIR). The large spatial coverage is necessary in order to obtain a large sample of clusters for carrying out a statistical analysis, and the broad spectral range allows accurate determinations of the individual cluster properties (Bik et al. 2003; Anders et al. anders04 (2004)). A preliminary analysis of a subset of the M51 cluster population was carried out by Bik et al. (2003; hereafter Paper I) who introduced the method used to determine the cluster properties and derived the age and mass distributions of the cluster sample roughly 2 kpc to the North East of the nucleus.
Bastian et al. (bastian05 (2005); hereafter Paper II) extended the survey to include the entire inner $`5`$ kpc of M51, and found 1152 clusters, 305 of which had accurate size determinations. In that work we extended the age distribution analysis of Paper I and found evidence for a cluster formation rate increase $`5070`$ Myr ago. This corresponds to the last close passage of NGC 5195 and M51 (Salo & Laurikainen salo00 (2000)). Additionally we found that 68 $`\pm 15`$% of the clusters forming in M51 will disrupt within the first $`10`$ Myr after their formation, independent of their mass, so-called infant mortality. For the resolved cluster sample, we found that the size distribution (the number of clusters as a function of their effective radius) can be well fit by a power-law: $`Ndr_{\mathrm{eff}}r_{\mathrm{eff}}^\eta dr_{\mathrm{eff}}`$, with $`\eta =2.2\pm 0.2`$, which is very similar to that found for Galactic globular clusters. Finally, we did not find any relation between the age and mass, mass and size, nor distance from the galactic center and cluster size.
In this study we focus on the evolution of the population of clusters in M51, in particular the timescale of cluster disruption and possible variations in the cluster formation rate. Cluster disruption of multi-aged populations, which excludes the galactic globular clusters, has been the subject of many earlier studies (e.g. Hodge (hodge87 (1987)) for the SMC and Battinelli & Capuzzo-Dolcetta (battinelli91 (1991)) for the Milky Way). In this study we will take into account mass dependent disruption, since the time needed to destroy half of the cluster population, which has been estimated in earlier work, will strongly depend on the mean mass of the sample and the lower mass limit of the sample. In addition, we here want to study the effect of variations in the formation rate, which is usually kept constant.
Boutloukos & Lamers (boutloukos03 (2003)) have developed a method to derive the disruption timescale based on the age and mass distributions of a magnitude limited cluster sample. They found that the disruption time of clusters in M51 is a factor of 15 shorter than the one for open clusters in the solar neighborhood. Lamers, Gieles & Portegies Zwart (lamers05a (2005)) showed that part of the difference can be explained by the difference in density of the cluster environment. They showed that the disruption time of clusters depends on the clusters initial mass and the galaxy density as $`t_{\mathrm{dis}}M_i^{0.62}\rho _{\mathrm{gal}}^{0.5}`$, based on the results of $`N`$-body simulations. The disruption time of clusters in M51 was still about a factor of 10 lower than the predicted value. In this work we are particularly interested if a short disruption timescale can be mimicked by an increasing cluster formation rate, and how the assumed disruption law influences the derived timescales. To this end we have generated artificial cluster samples with parameterized global characteristics (e.g. time dependent cluster formation rates, disruption laws, infant mortality rates, and mass functions). We then compare these models with the derived age and mass distributions of the cluster population of M51 to derive the best fit parameters for the population as a whole.
The structure of the paper is as follows: In § 2 the observations of the cluster population of M51 are presented. In § 3 we investigate how the disruption time depends on different cluster and galaxy parameters. § 4 describes the steps we will take in our models, where the details of the models we used to generate artificial cluster populations will be explained in § 5. The results of the fits are given in § 6. A discussion on the implication of the results is given in § 7. The conclusions are presented in § 8.
## 2 The observations
### 2.1 Fitting the observed spectral energy distribution
From archival HST broadband photometry we have derived the age, mass, and extinction of 1152 clusters in M51 (Paper II) using the three-dimensional maximum likelihood fitting (3DEF) method. Details about the 3DEF method can be found in Paper I. In summary, the spectral energy distribution of each cluster is compared with cluster evolution models. In this case the GALEV simple stellar population (SSP) models (Anders et al. anders03 (2003); Schulz et al. schulz02 (2002)) for solar metallicity and Salpeter IMF are used. For each age a series of different extinctions is then applied to the models and all combinations of age and extinction are compared to the data. The lowest $`\chi ^2`$ is kept and from the absolute magnitude at that age the mass is determined. Detailed tests of the accuracy and reliability of the derived parameters are presented in Paper II. In the present work we further investigate the accuracy of our fitting method and we use the data set to develop a model that describes the global properties of the cluster system.
### 2.2 The age and mass distribution of clusters
The ages and present masses of the 1152 clusters are plotted in the top panel of Fig. 1. In order to be able to compare our observations with simulated cluster samples, we bin the data in logarithmic number density plots of the age vs. mass distribution. Clusters where counted in bins of 0.4 age dex by 0.4 mass dex (Fig. 1, Middle). The result is illustrated in the bottom panel of Fig. 1. A few striking features can be learned from this diagram:
1. a burst in the cluster formation rate (CFR) between 50 and 70 Myr, corresponding with the most recent interaction with the companion galaxy NGC 5195;
2. a short lived young population with ages $`<`$ 10 Myr. In Paper II we found that $`\pm `$68% of these young clusters will dissolve within 10 Myr, independent of their mass;
3. evolutionary fading under the detection limit, which makes it harder to detect old low mass clusters. The increasing line in Fig. 1 shows how the 90% completeness limit in the F439W band (22.6 mag.) corresponds with different masses at different ages;
4. an apparent increase in the mass of the most massive cluster with age. This is a binning effect: the older age bins span more time and therefore contain more clusters and the chance of finding a more massive cluster at older ages is higher due to the size of sample effect (Hunter et al. hunter03 (2003)). We have tried to use the method of Hunter et al. (hunter03 (2003)) to derive the cluster formation rate, but we found that the increase in the maximum mass is much to shallow. This probably means that M51 has reached the maximum cluster mass ($`10^6\text{M}_{}`$) and then the relation of Hunter et al. (hunter03 (2003)) does not apply anymore, since the maximum cluster mass found at a certain age is then not determined anymore by sampling statistics. This is the topic of a next study (Gieles et al. gieles05 (2005)).
We found that a large fraction ($`\pm 68`$%) of the clusters younger than $`10^7`$ yr dissolves independent of mass, probably due to the removal of the primordial gas (e.g. Kroupa kroupa04 (2004); Geyer & Burkert geyer01 (2001); Lada & Lada lada03 (2003)). After these critical 10<sup>7</sup> years, the surviving clusters will dissolve due to the tidal field of the host galaxy and external perturbations from, for example, encounters with giant molecular clouds (GMCs). The disruption of clusters is in that sense a two step process. Here we will use the resulting age/mass distribution to study the disruption of clusters with ages larger than $`10^7`$ yr., i.e. the second step in the disruption process.
### 2.3 Artifacts introduced by the age fitting method
We want to see whether the 3DEF method (§ 2.1), used to derive ages, masses and extinctions from the photometry introduces systematic artifacts. More important, could it affect our results of the disruption time or formation rate? For instance, are there systematically old clusters fitted with young ages or the other way around?
The uncertainty in the derived ages, extinctions and masses from broad-band photometry is mainly caused by two effects:
1. Differences between the real integrated colors of the clusters and the models, for example due to stochastic sampling of the stellar IMF which cannot be taken into account in the SSP models, or errors in the stellar isochrones used;
2. Systematic errors introduced by the age fitting method and the applied selection effects
Of course the first effect cannot be corrected for, unless we are able to compare photometric determined ages directly with spectroscopically determined ages, which is unfortunately not feasible for a whole population of clusters. In addition, although spectroscopically derived ages are more accurate, they are hampered by their own problems (Brodie et al. brodie98 (1998)). In addition, the age derivation will dependent on the choice of the adopted SSP models. Variations in the metallicity and the IMF will affect the derived ages. A detailed comparison of the data with models of different metallicity is carried out in Paper II. An earlier study of the age distribution of M51 used Starburst99 models (Bastian & Lamers bastian03 (2003)), and they found a very similar age distribution.
The second effect can be quantified with the use of artificial cluster populations. Earlier studies (e.g. Anders et al. anders04 (2004); de Grijs et al. degrijs05 (2005)) have already shown the importance of using a long wavelength baseline ($`U`$ to $`NIR`$) for age dating young clusters. Here we will make an attempt to quantify possible systematic errors introduced by the age-fitting method and see whether we can correct for them or not.
To quantify the artifacts introduced by the fitting routine, an artificial cluster sample including simulated observational errors and extinction values is generated and fitted with the same fitting procedure as used for the data (§ 2.1). We start with a sample of clusters equally spread in log(Age/yr) and log($`M/\text{M}_{}`$) space. In total 201 time steps between log(Age/yr) = 6 and log(Age/yr) = 10 and 161 mass-steps between log($`M/\text{M}_{}`$) = 2 and log($`M/\text{M}_{}`$) = 7 were generated. The GALEV models have log(Age/yr) = 6.6 as youngest model, so clusters with younger ages were given that age. The magnitudes as a function of age and mass where taken from the GALEV SSP models. Observational uncertainties were applied as a function of magnitude as was done in Paper II: the observed errors in the magnitudes of clusters in M51 can be well approximated by $`\mathrm{\Delta }\mathrm{mag}_\lambda =10^{d_1+d_2\times \mathrm{mag}_\lambda }`$. The values for $`d_1`$ and $`d_2`$ for the different filters are results from analytical functions fitted to the observed errors and magnitudes and are given in Table 4 of Paper II. Ideally, we then apply the same extinction to the model clusters as the M51 clusters have. Unfortunately the only information we have is the extinction we have measured, which of course could already be polluted with artifacts. To get an estimate of the uncertainty in the measured extinction, we start with a sample of clusters with no extinction applied. When we fit this population with the 3DEF method, we find that 20% of the sources is fitted with some extinction. This is quite a large number, but fortunately 90% of these sources have extinction values lower than $`E(BV)`$ = 0.1 mag. The maximum $`E(BV)`$ found is 1 mag.
The next step is to apply an extinction model close to what we observe. To this end $`E(BV)`$ extinction values were chosen randomly from a Gaussian distribution centered at 0 with $`\sigma =0.10`$ for clusters younger than log(Age/yr) = 7.3 and $`\sigma =0.05`$ for clusters older than log(Age/yr) = 7.3. The values for $`\sigma `$ agree with the value we found for the mean extinction in Paper II. There we found that these values are the average extinction for these two age groups. The higher extinction for young ages is caused by the presence of the left-over dust around the cluster. Negative extinctions were set to 0, resembling the extinction distribution of the data where half of the clusters had $`E(BV)`$ = 0 (See Fig. 8 of Paper II). An age dependent maximum extinction was applied of the form: $`E(BV)_{\mathrm{max}}(t)=50.5\mathrm{log}(t)`$. This is a little bit lower than the observed maximum extinction, but we know that some of the observed high values could be caused by wrong fits. This still resembles the observed extinction behavior quite well. The resulting magnitudes are than cut off at our completeness limits in each filter. In this way we created the spectral energy distributions of a large artificial cluster sample with age, mass and extinction known for each cluster. These were fitted with the 3DEF method (§ 2.1).
The result of the fitted simulation is shown in Fig. 2. A direct comparison with the observed age-mass diagram of M51 clusters (Fig. 1, Top) shows that there are features present in the data which are not visible in the fitted simulation. For example, there is a gap at 6.9 $`<`$ log(Age/yr) $`<`$ 7.1 in the M51 cluster sample which seems to appear at slightly higher ages in the fitted simulations (7.1 $`<`$ log(Age/yr) $`<`$ 7.2 ). This suggests that the artifacts in the data are not only caused by our applied selection effects or our age-fitting technique. In the top left panel of Fig. 3 we show the fitted age versus the input age for the simulated cluster sample. A large number of clusters with wrong ages are fitted with an log(Age/yr) $``$ 7. We found for 87% of the modeled clusters that the fitted age was the same as the input age within 0.4 dex (Fig. 3, Top right). For the mass 97% was fitted correctly within 0.4 dex (Fig. 3, Bottom right) and 92% of the extinction values where fitted back within 0.05 mag. (Fig. 3, Bottom left). We have to realize that the strength of the artifacts depend on the number of input clusters at each age and mass bin. We have not attempted to match the observations in this stage, since we are only interested in relative errors. For example, the horizontal spur at log(Age/yr) $``$ 7 in Fig. 3 (Top, left) is populated with clusters with input ages up till a few times 10<sup>9</sup> yr. The number of clusters with that age in our M51 sample is very low (See Fig. 1, Top).
### 2.4 Correcting the data for fitting artifacts
We will try to correct the data for possible artifacts, by using the (systematic) deviations found in § 2.3. Correcting the observed ages based on the absolute numbers deviating from the one-to-one relation is not useful, since the number of clusters that were used as input at each age and mass differs from the observed number. From the input sample we can derive how many clusters are (systematically) fitted with wrong ages and masses. Let the total number of bins in age/mass space be $`K`$. Here we define the number of bins as the number of bins which are not affected by the detection limit (See Fig. 1, Middle panel). The number of clusters found in each bin is the sum of the contribution of clusters from all bins to this one, where the majority will be from the bin with the same input age and mass. When we write the number of clusters in each bin as a vector with $`K`$-entries, the fitted number of clusters can be written as a matrix multiplication of all contributions times the input number of clusters
$$\left(\begin{array}{c}N_1\\ N_2\\ \mathrm{}\\ N_K\end{array}\right)_{\mathrm{obs}}=\left(\begin{array}{cccc}C_{11}& C_{12}& \mathrm{}& C_{1K}\\ C_{21}& C_{22}& \mathrm{}& C_{2K}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ C_{K1}& C_{K2}& \mathrm{}& C_{KK}\end{array}\right)\left(\begin{array}{c}N_1\\ N_2\\ \mathrm{}\\ N_K\end{array}\right)_{\mathrm{intr}}$$
(1)
where $`N_{\mathrm{intr},j}`$ is the number of clusters generated in bin $`j`$, $`N_{\mathrm{obs},i}`$ is the number of clusters fitted in bin $`i`$ and $`C_{ij}`$ is the contribution to bin $`i`$ from bin $`j`$. Tests have shown that the best results are acquired when only taking bins into account which are not affected by the detection limit (see Fig. 1, Middle panel). From the simulated and fitted sample we can derive the values for $`C_{ij}`$ for all combinations of $`i`$ and $`j`$. All values on the diagonal of $`C`$ (i.e. $`C_{ij}`$ where $`i`$ = $`j`$) are close to 1 since most clusters are fitted with the same age, extinction and mass. All other values are 0 or between 0 and 1. When we know the matrix $`C`$, the inverse can be used to correct the observations for systematic 3DEF fitting artifacts.
$$N_{\mathrm{intr}}=C^1\times N_{\mathrm{obs}}$$
(2)
where $`N_{\mathrm{intr}}`$ is the vector with the intrinsic number of clusters, $`C^1`$ is the inverse of the contribution matrix as defined in Eq. 1 and $`N_{\mathrm{obs}}`$ is the vector with observed clusters. When we correct the observed vector and plot it in a 2D age-mass diagram again, divide the corrected and uncorrected observation, we can see where deviations take place. Fig. 4 shows the ratio of corrected over uncorrected observations. The number of clusters in the age bins at log(Age/yr) = 7.2 and 8.0 has been lowered with about 15%. The observations, however, show a gap at log(Age/yr) = 7.2 (Fig. 1, Top). So, the underestimation of clusters in that age bin is not caused by our fitting routine.
The corrected observations based on Eq. 2 are shown in Fig. 5. The burst at between 50-70 Myr is less pronounced, but still present. The differences with the uncorrected observations (Fig. 1, Bottom) are small, and we therefore conclude that our age-fitting method (3DEF) and our applied selection affect is not severely affecting our age-mass diagrams in a systematic way. Especially, there is no large systematic shift from old to young clusters or the other way around. We therefore conclude that we can use the uncorrected data as well as the corrected data to compare with the synthetic cluster populations in § 5. In § 6.2 we will show that both the corrected as the uncorrected data give the same results when fitting the analytical models to the data.
## 3 Exploration of the parameters which determine the disruption time
### 3.1 Cluster initial mass
If cluster relaxation drives the evaporation of clusters, then the more massive clusters live longer than their low mass counterparts. Boutloukos & Lamers (boutloukos03 (2003)) have proposed an empirical way to determine the dependence of the cluster disruption time on the initial cluster mass, assuming a power-law dependence of the disruption on the cluster mass
$$t_{\mathrm{dis}}=t_4(M_i/10^4\text{M}_{})^\gamma $$
(3)
where $`t_4`$ is the disruption time of a $`10^4\text{M}_{}`$ cluster, $`M_i`$ is the initial mass of the cluster and $`\gamma `$ is a dimensionless index. The value of $`\gamma `$ was determined by measuring the slope of the age and mass distributions of a cluster population. Their mean value for $`\gamma `$ based on four different galaxies was $`<\gamma >=0.62\pm 0.06`$. In a recent study by Lamers, Gieles & Portegies Zwart (lamers05a (2005)) these observational results are compared with results of $`N`$-body simulations. The value of $`\gamma `$ can be explained by tidally driven relaxation and was confirmed to be 0.62 by $`N`$-body simulations. The explanation for this is that the disruption time of a cluster in a tidal field, depends on the relaxation time ($`t_{\mathrm{rel}}`$) and the crossing time ($`t_{\mathrm{cr}}`$) of the cluster as (Baumgardt baumgardt01 (2001)):
$$t_{\mathrm{dis}}t_{\mathrm{rel}}^x\times t_{\mathrm{cr}}^{1x}$$
(4)
Using the expression for $`t_{\mathrm{rel}}`$ and $`t_{\mathrm{cr}}`$ from Spitzer (spitzer87 (1987)), this implies: $`t_{\mathrm{dis}}\beta (N/\mathrm{ln}N)^x`$. Baumgardt & Makino (baumgardt03 (2003)) found two combinations of $`\beta `$ and $`x`$, depending on the concentration of the clusters. Lamers, Gieles & Portegies Zwart (lamers05a (2005)) showed that for both combinations $`t_{\mathrm{dis}}`$ could be well approximated with $`t_{\mathrm{dis}}N^{0.62}M^{0.62}`$, where $`N`$ is the number of stars in the clusters and $`M`$ is the total mass of the cluster. King (king58 (1958)) already had theoretical arguments to expect that the lifetime of clusters should depend on the mass as $`t_{\mathrm{dis}}M^{2/3}`$. The agreement between observations and $`N`$-body simulations was first noted by Gieles et al. (gieles04 (2004)). The physical background of why the disruption time does not scale directly with the relaxation time is given by Fukushige & Heggie (fukushige00 (2000)).
### 3.2 Cluster radius
Young clusters are not only affected by the external tidal field of the host galaxy, but they also undergo shocks from (giant) molecular clouds. Both these processes shorten the lifetime of clusters. For both cases the radius of the cluster is an important parameter in determining how fast the cluster will disrupt. However, both processes depend very different on the radius. From Eq. 4 and the expression for the relaxation time and crossing time it follows that for the tidally driven relaxation the disruption time depends on the radius as $`t_{\mathrm{dis}}r_\mathrm{h}^{3/2}`$. Larger clusters live longer since they have a longer relaxation time, so it takes more time for stars to reach the tidal radius and leave the cluster. Spitzer (spitzer58 (1958)) has shown that the time needed for a cluster to get unbound due to external shocks, relates to the half mass radius of the cluster as: $`t_{\mathrm{sh}}r_\mathrm{h}^3`$, so here larger clusters live shorter (for isolated clusters).
To see whether the radius of a cluster is an important parameter in disruption, we use the radii measurements of Paper II. There we measured the projected half light radius (or effective radius) $`r_{\mathrm{eff}}`$, which relates to the half mass radius as: $`r_{\mathrm{eff}}=3/4r_\mathrm{h}`$ (Spitzer spitzer87 (1987)). We make a number density plot of $`r_{\mathrm{eff}}`$ vs. age for all clusters (Fig. 6). There are clearly no old clusters with large radius, while the opposite is expected due to the size-of-sample effect (Hunter et al. hunter03 (2003)). This suggests that large clusters are disrupted preferentially. This suggests that shocks may be the dominating disruption effect. However, when a large fraction of the clusters is removed, independent of radius, the upper radius would also go down. This is a result from number statistics: less clusters in a power law distribution will result in a lower maximum value. So what really matters here is whether the slope of the radius distribution changes in time or not. In Paper II it is shown that the cluster radius distribution of M51 is $`N(r)drr^\eta dr`$, with $`\eta =2.2\pm 0.2`$. To see how the slope of the distribution depends on age, we divide our cluster sample in young (log(Age/yr) $`<`$ 7.5) and old (log(Age/yr) $`>`$ 7.5). Dividing the sample at log(Age/yr) = 7.5 yields two samples of more or less equal size, which gives similar errros in the fit to both distributions. When we determine this index $`\eta `$ for only young clusters we find $`\eta =2.0\pm 0.4`$ and for old clusters we find $`\eta =2.5\pm 0.6`$, which is very similar for the value found for the globular clusters in our Milky Way ($`\eta =2.4\pm 0.2`$, Paper II). Although the radius distribution seems to get steeper with age, the errors are too large to place a strong constraint on this. We therefore do not take the radius into account as a free parameter when modeling the cluster disruption. Futher studies of M51 with higher resolution, for example with the Advanced Camera for Surveys (ACS), could shed light on how the radius of clusters affects the lifetime.
### 3.3 Distance to the galactic center
Lamers, Gieles & Portegies Zwart (lamers05a (2005)) and Baumgardt & Makino (baumgardt03 (2003)) have shown that the disruption time is expected to depend on the galactocentric distance of the cluster, the orbital velocity in the galaxy and the ambient density of the galaxy as
$`t_{\mathrm{dis}}`$ $``$ $`R_\mathrm{G}/V`$ (5)
$``$ $`\rho _{\mathrm{amb}}^{0.5}`$ (6)
where $`R_\mathrm{G}`$ is the distance to the galactic center, $`V`$ is the rotational velocity of the cluster in the host galaxy at that distance and $`\rho _{\mathrm{amb}}`$ is the ambient density of the galaxy at the location of the cluster. The relation with the ambient density holds only when a logarithmic potential is assumed. Since we are dealing with a disk galaxy, it is not so straightforward to derive the ambient density from the spherically symmetric logarithmic potential. Therefore, we prefer the relation with the galactocentric distance and the velocity (Eq. 5). With Eq. 5 we are able to estimate if we would be able to observe a difference in disruption time at different locations in the galaxy. When we look at clusters between 1 and 3 kpc and at clusters between 3 and 5 kpc, the average value of $`R_G`$ goes up with a factor of 2. The rotational velocity of M51 increases from 200 km s<sup>-1</sup> to 225 km s<sup>-1</sup> (Rand rand93 (1993)). So from Eq. 5 we expect the disruption time in the two samples to be different by a factor of 1.8. In Fig. 7 we plot the ratio of the number of clusters at different ages for the outer region (3-5 kpc) and inner region (1-3 kpc). Overplotted is the predicted ratio using Eq. 8 for disruption timescale that is different by a factor of 1.8. Just as for the radius dependence we see that the distance to the galactic center plays a role, but the data is not sufficient to include it in our analysis.
In conclusion, we see evidence for radius and galactocentric distance dependent disruption, but the noise is too large to include these parameters in the models. The mass of the cluster is the most dominating parameter in the determination of disruption time and in the remainder of the study we will only use the mass dependence as a parameter we will vary in the models.
## 4 Input parameters for modeling the cluster population of M51
So far, analytical models for finding the cluster disruption time have assumed that clusters were formed with a constant CFR, as is probably the case for Galactic open clusters (Boutloukos & Lamers boutloukos03 (2003); Battinelli & Capuzzo-Dolcetta battinelli91 (1991)). Lamers et al. (lamers05b (2005)) predicted the age distribution of open clusters. In the case of M51, we have age and mass information available for each cluster, so predictions can be done for age and mass. In addition, assuming a constant CFR for M51 might be an oversimplification of the situation, since the galaxy is in interaction with NGC 5195. In the next sections we explore a broader parameter space.
Since there are strong arguments to believe that the mass dependence of the cluster disruption ($`\gamma =0.62`$) is constant (Lamers, Gieles & Portegies Zwart lamers05a (2005)), we start by varying only the constant $`t_4`$, to be able to compare our results with clusters gradually loosing mass with the instantaneous disruption assumption (Eq. 3) results of Boutloukos & Lamers (boutloukos03 (2003)). Next, a two dimensional parameter search for $`\gamma `$ and $`t_4`$ is performed, to verify the assumed value for $`\gamma `$ and to study the dependence of $`t_4`$ on the value of $`\gamma `$.
When we have a first estimate of the disruption time, we will study how this value changes when we assume that the CFR has been increasing during the last Gyr or contains bursts at the moments of encounter with NGC 5195.
## 5 Analytical model for generating a cluster population
### 5.1 Setting up a synthetic cluster population
The synthetic cluster populations will be created in a similar way as in § 2.3. This time however we want to include realistic input physics, like the cluster IMF and different formation rates, so creating clusters equally spaced in log(Age/yr) and log($`M/\text{M}_{}`$) will not be adequate. When creating clusters with a realistic CIMF, the number of clusters needed to fully sample the CIMF up to 10 Gyr ago is too high. Therefore each cluster was assigned a weight depending on the initial cluster mass and its age ($`w(t,M_i)`$), proportional to the expected number of clusters formed at each age and mass. The weight is a function of age and mass, scaled such that the youngest most massive cluster has a weight of 1
$$w(t,M_i)=(t/t_{\mathrm{min}})\times (\alpha 1)\times (M_i/M_{\mathrm{max}})^{1\alpha }$$
(7)
where $`w(t,M_i)`$ is the weight assigned to a cluster with age $`t`$ and mass $`M_i`$, $`M_{\mathrm{max}}`$ is the mass of the most massive cluster in the simulation, $`t_{\mathrm{min}}`$ is the age of the youngest cluster in the simulation and $`\alpha `$ is the slope of the mass function. When $`\alpha `$ is chosen 2, i.e. $`N(M)M^2`$, the weight depends on age and mass simply as: $`w(t,M_i)t/M_i`$. When the simulated clusters are binned, the weights of the clusters are counted, yielding a realistic log(Age/yr) vs. log($`M/\text{M}_{}`$) diagram similar to Fig. 1 (Bottom). The advantage of using points spread equally in log(Age/yr) and log($`M/\text{M}_{}`$) with weights assigned, is that the number of points per bin is constant and that it is very easy to create a lot of populations with different formation rates, disruption timescales etc. in a short time.
In our case the clusters are given a weight such that after binning the CIMF has a slope of $`\alpha =2.1`$ as found for M51 (Paper I) and the Galactic open clusters (Battinelli et al. battinelli94 (1994)) and using different formation and disruption scenarios (§ 6.2-§ 6.4).
### 5.2 Including stellar evolution and cluster disruption
Baumgardt & Makino (baumgardt03 (2003)) have shown that stellar evolution (SEV) is an important contributor to the dissolution of young clusters, especially for clusters with low concentrations. They also confirmed that clusters dissolve with a power-law dependence of their initial mass as $`t_{\mathrm{dis}}M_i^\gamma `$, where $`\gamma =0.62`$, in agreement with the empirical determination by Boutloukos & Lamers (boutloukos03 (2003)). In the latter study instantaneous disruption after the disruption time was assumed as a first approximation and they found that the typical disruption time ($`t_4`$, see Eq. 3) varies a lot for different galaxies. In a recent study (Lamers et al. lamers05b (2005)) it was shown that there is a simple analytical description of the mass of a cluster as a function of time. It takes into account the effect of mass loss due to stellar evolution, based on the mass loss predicted by the GALEV SSP models (Anders et al. anders03 (2003); Schulz et al. schulz02 (2002)) and cluster mass loss due to the tidal fields. The mass of the cluster as a function of time can be well approximated by
$$M_p(t)=((M_i\mu _{\mathrm{sev}})^\gamma \gamma \frac{t}{t_0})^{1/\gamma }$$
(8)
where $`M_p(t)`$ is the present mass of the cluster as a function of its age, $`M_i`$ is the initial mass of the cluster, $`\mu _{\mathrm{sev}}M_p(t)/M_i`$ is the fraction of remaining mass after mass loss due to stellar evolution has been subtracted and $`t_0`$ relates to $`t_4`$ as $`t_4=t_0\times 10^{4\gamma }`$. The mass as a function of time, according to the analytical formula, agrees perfectly with the predictions following from $`N`$-body simulations. Lamers et al. (lamers05b (2005)) have also shown that with this analytical model the age distribution of galactic open clusters can be explained very well.
## 6 Fitting observed age-mass distribution to predictions
### 6.1 Determining reduced $`\chi ^2`$ values from 2D fits
Artificial cluster samples with realistic input physics (e.g. a CIMF, cluster disruption, bursts etc.) can now be generated and compared with the observed age and mass number density distribution.
After calculating the analytically generated cluster population, the model is binned into number density plots in the same way as the observed data (see § 2.2) taking into account the weights. In order to compare the simulated (2D) age-mass density plots with the observations, we use the Poisson Probability Law (PPL) introduced by Dolphin & Kennicutt (dolphin02 (2002)) for similar purposes
$$\mathrm{PPL}=2\underset{i=0}{\overset{N}{}}m_in_i+n_i\mathrm{ln}\frac{n_i}{m_i}$$
(9)
where $`N`$ is the number of bins, $`m_i`$ is the predicted number by the analytical model in bin $`i`$ and $`n_i`$ is the observed number of clusters in bin $`i`$. The value of PPL is similar to the $`\chi ^2`$, in the sense that lower values imply better fits. We will always divide the PPL value by the number of bins minus the number of degrees of freedom, which is equivalent to the reduced $`\chi ^2`$, so the $`\chi _\nu ^2`$. We will refer to $`\chi _\nu ^2`$ when we discuss results of fits.
### 6.2 Determining the cluster disruption time assuming a constant formation rate of clusters
To determine the typical cluster disruption time, $`t_4`$, defined in § 5, we generate a cluster sample with a constant CFR and then calculate the cluster masses as a function of age according to Eq. 8 for various values of $`t_4`$. Here we are interested in the disruption time of clusters that have survived that first 10<sup>7</sup> yr in which the natal cloud is being removed by stellar winds, therefore we exclude the youngest age bin in the fits. Fig. 8 shows a clear $`\chi _\nu ^2`$ minimum around $`t_4=1.0_{0.5}^{+0.6}\times 10^8`$ yr, where the upper and lower errors are defined by $`\chi _{\nu }^{}{}_{,\mathrm{accept}}{}^{2}\chi _{\nu }^{}{}_{,\mathrm{min}}{}^{2}+1`$, which is equivalent to the 1 $`\sigma `$ error. In addition, we have fitted the same models but then corrected for age-fitting artefacts (§ 2.4). The shape of this $`\chi _\nu ^2`$ curve is the same as for the raw data, though the values are higher. This shows that the uncertainties of our age-fitting method do not alter the value found for the disruption timescale.
To see how the value of $`t_4`$ depends on the value of $`\gamma `$, we simulate a grid of cluster populations and vary $`t_4`$ and $`\gamma `$. A 2D $`\chi _\nu ^2`$ plot is shown in Fig. 9. The minimum is at $`\gamma =0.65_{0.25}^{+0.16}`$ and $`t_4=1.0_{0.35}^{+0.84}\times 10^8`$ yr, agreeing very well with the value of $`\gamma =0.62`$, which was stated earlier based on theoretical arguments and other observational results. The plot also shows that there is a diagonal bar-shaped minimum for different combinations of $`t_4`$ and $`\gamma `$. One could argue that there could be multiple combinations possible, which will yield a somewhat higher value for $`t_4`$. The fit however is very sensitive for the choice of bin size when varying two variables. We excluded the mass bins higher than $`5\times 10^5\text{M}_{}`$, since we probably deal with a truncation of the mass function. If sampling effects would determine the upper mass at different ages (Hunter et al. hunter03 (2003)), the maximum mass should increase much more than we observe in the top panel of Fig. 1. This effect makes the mass function steeper above log($`M/\text{M}_{}`$) $`5.3`$ and therefor that region in the age/mass diagram is not suitable to fit the (sensitive) mass dependent disruption. An alternative way to measure $`\gamma `$ would be to measure the slopes of the age and mass distribution separately, as was done in Boutloukos & Lamers (boutloukos03 (2003)). We fitted these slopes and found the same value for $`\gamma `$ as for the 2D fit shown in Fig. 9. Again, for the mass, we do not include the high mass end for similar reasons as mentioned before. This method is less sensitive for the choice of bin size, since we can fit the slope of the age and mass distribution independent of the value of the disruption time. We choose to include the result of the simultaneous fit of $`t_4`$ and $`\gamma `$, because it illustrated nicely how these two variables relate.
### 6.3 The effect of an increasing cluster formation rate
Since NGC 5195 is probably bound to M51 and therefore slowly falling in (Salo & Laurikainen salo00 (2000)), one could argue that the short disruption timescale found in § 6.2 is actually caused by an increasing cluster formation rate (CFR). Bergvall et al. (bergvall03 (2003)) have shown that interacting galaxies such as M51 (i.e. non-merging), can have an increased star formation rate of the order of a factor of 2-3. We therefore model different cluster populations with increasing CFR($`t`$) rates of various strengths, where we assume that an increasing star formation rate results in an equally large increase in the CFR($`t`$). We study two different models with increasing CFR: 1.) a linear increasing CFR starting 1 Gyr ago (§ 6.3.1); 2.) a CFR that increases with bursts at the moments of encounter with NGC 5195 (§ 6.3.2). Fig. 10 gives a schematic illustration of how the CFR varies with time for the two models.
#### 6.3.1 Linearly increasing cluster formation rate
In the linear model the CFR($`t`$) starts to increase 1 Gyr ago, which is before the moment of the early close encounter with NGC 5195 (400-500 Myr ago, Salo & Laurikainen salo00 (2000)). We expect the CFR to start increasing before the moment of the closest encounter, since the two galaxies are already interacting before the first perigalactic passage. We calculate models with different CFR increases and disruption times. We plot the $`\chi _\nu ^2`$ values for various values of CFR($`t`$ = 0)/CFR($`t`$ = 10<sup>9</sup> yr) and $`t_4`$ in top panel of Fig. 11.
The minimum $`\chi _\nu ^2`$ value is at $`t_4=2.0_{1.1}^{+5.2}\times 10^8`$ yr and an increase in the CFR of $`7.0_{5.0}^{+68.1}`$. For small values of $`t_4`$ the equal $`\chi _\nu ^2`$ lines are vertical. This can be explained by the fact that if the disruption time is short, no fingerprints of the ancient formation rate are present in the current population. They are simply erased by disruption. the reason that the equal $`\chi _\nu ^2`$ contours are circular around the minimum, is that the disruption of clusters depends on the mass of the clusters (Eq. 3), unlike an increase of the formation rate.
#### 6.3.2 Cluster formation rate with bursts
An alternative formation scenario would be that the CFR increases with a burst at the moments of encounter with NGC 5195 and then an exponential decay in the CFR (see model 2 in Fig. 10). We choose the moments of increase at $`t=7\times 10^7`$ yr and $`t=5\times 10^8`$ yr ago, based on the results of Salo & Laurikainen (salo00 (2000)) and the typical decay time of the burst is $`10^8`$ yr (Paper II). The CFR step and $`t_4`$ are varied in different models. The bottom panel of Fig. 11 shows that the lowest $`\chi _\nu ^2`$ value is at $`t_4=2.0_{1.1}^{+2.3}\times 10^8`$ yr. This is a factor of 2 higher than when the increasing CFR is not taken into account, but it is the same value as was found for the linear increase in the CFR. The value is still a factor of 5 lower then predicted by $`N`$-body simulations (Baumgardt & Makino baumgardt03 (2003); Lamers, Gieles & Portegies Zwart lamers05a (2005)). The best value for the increase in CFR at the moment of encounter is $`3.0_{1.2}^{+4.6}`$. The latter value agrees very well with what is generally observed for the increase in star formation rate of interacting galaxies (Bergvall et al. bergvall03 (2003)). Since one of the bursts is clearly observed and a linearly increasing CFR is not so physical, we prefer Model 2 above Model 1. In the next section we will compare several properties of this model with the observations.
### 6.4 Comparision between the best fit model and the observations
We show a direct comparison between the age-mass diagrams of the best fit model (§ 6.3.2) and the observations in Fig. 12. The densities are scaled such that the total number of simulated clusters equals the total number of observed clusters (1152). A few bins in the observations are empty and not empty in the simulations. The reason for this is that the simulated cluster sample containts bins with values smaller than 1. Apart from this, the general trend of grey values in this 2D plot is very similar in both cases.
Another interesting property of the observations is the formation rate. In Paper II we showed the number of clusters at different ages for different mass cut offs. For clusters with masses higher than $`10^{4.7}\text{M}_{}`$ we get a realistic impression of the cluster formation rate. This is because we are complete until 1 Gyr for these masses (see top panel of Fig. 1) and because the most massive clusters are not affected by disruption that much. In Fig. 13 (Top) we show the number of clusters in different age bins for the observations and the best fit model. The general trend of the observations is followed very well by the model. A better way to show the formation rate is to divide each age bin by the width of the bin. Then we get the number of clusters formed per unit of time (Myr). This is shown in the bottom panel of Fig. 13. In this figure the over-density of young clusters (log(Age/yr) $`<`$ 7) is more obvious and the burst at $`7\times 10^7`$ year is better visible. The first burst of cluster formation ($`5\times 10^8`$ years ago) is not visible anymore, since clusters with these ages are already affected by the (short) disruption time. This reinforces that it is very hard to detect variations in the cluster formation rate when the disruption time is that short. The largest difference is seen for the bin with log(Age/yr) = 7 and 8.25. The model predicts in these bins more clusters than are observed. This can be explained by fitting artefacts which yield an (unphysical) underdensity of clusters (§ 2.3). The model is still within the 3 $`\sigma `$ error of the observations however.
## 7 Implication of the derived disruption time
We have shown that the clusters disruption time for a typical cluster with mass of $`10^4\text{M}_{}`$ is around $`10^8`$ years in M51. When increasing formation rates are taken into account the disruption time increases with a factor of 2. This is significantly longer than Boutloukos & Lamers (boutloukos03 (2003)) found for clusters in a smaller region of M51 ($`t_4=4\times 10^7`$ yr). This could be because they did not seperate the dissolution due to infant mortality rate from the evaporation by the tidal field from the galaxy. Clusters with ages younger than $`10^7`$ year are not taken into account in this study, since they are affected by the dissolution due to the removal of primordial gas. Theoretical predictions show that clusters of $`10^4\text{M}_{}`$ in a tidal field of the strength of M51 should have an disruption time of about $`10^9`$ years (Lamers, Gieles & Portegies Zwart lamers05a (2005); Baumgardt & Makino baumgardt03 (2003)). This value is found from observations of clusters in the solar neighborhood (Boutloukus & Lamers boutloukos03 (2003)). What causes the clusters in M51 to dissolve about 5 times faster than predicted? A few effects that have not been incorporated in the $`N`$-body models that predict the disruption times in tidal fields are:
1. Variations in the stellar IMF. When clusters are formed with a so-called top heavy IMF as is observed in the starburst galaxy M82 (Smith & Gallagher smith01 (2001)), clusters will disperse much faster, since the disruption time depends on the number of stars in the clusters as $`t_{\mathrm{dis}}N^{0.62}`$ (§ 3.1). Suppose the stellar IMF starts at 1 $`\text{M}_{}`$ in stead of $`0.1\text{M}_{}`$, then the number of stars for a given cluster mass will be about a factor 10 lower. This will make the disruption time a factor $`10^{0.62}4`$ lower. This would nicely explain the factor 5 difference in disruption time we observe.
2. External perturbations. The $`N`$-body models of Baumgardt & Makino (baumgardt03 (2003)) calculated the disruption time of clusters in a smooth external potential from the host galaxy. In reality, the cluster will also experience additional external perturbations, for example the encounters with molecular clouds. The clusters in our sample are in the inner 5 kpc of the galaxy, where most of the giant molecular clouds reside (Henry et al. henry03 (2003); Kuno et al. kuno95 (1995)). The encounters with molecular clouds can speed up the disruption of clouds significantly (e.g. Terlevich terlevich87 (1987); Theuns theuns91 (1991)).
3. Out of equilibrium formation of clusters. All clusters in the $`N`$-body models start in virial equilibrium and in tidal equilibrium with the host galaxy. Kroupa (kroupa04 (2004)) has shown that after the gas removal phase, the clusters are not in virial equilibrium anymore and the outer parts of the cluster have expanded. It will be easier to dissolve these clusters than when all stars are in tidal equilibrium and within the tidal radius imposed by the host galaxy.
4. Variations in the central concentration. The $`N`$-body models of Baumgardt & Makino (baumgardt03 (2003)) start clusters with concentration values of $`W_0`$ = 5-7. This is the average concentration of globular clusters in our Milky Way (Harris harris96 (1996)). When clusters start with much smaller concentration, the core of the clusters is less compact and the cluster will be more vulnerable for external perturbations. The concentration of clusters in M51 can not be determined due to lack of resolution, but we know from young open clusters in the Milky Way that they have much smaller concentration indices than the globular clusters (Binney & Tremaine binney87 (1987)).
So far the clusters in M51 have not been checked for variations in the IMF in the way that it has been done for clusters in other galaxies (e.g. Smith & Gallagher smith01 (2001); Larsen et al. larsen04 (2004); Maraston et al. maraston04 (2004)). Also, no $`N`$-body experiments have been performed including the effect of a tidal field and perturbations by giant molecular clouds. Argument 3 and 4 are based upon unknown observables of young clusters and they could hold for clusters in other galaxies as well.
When the disruption of clusters is indeed as short as we derived, young massive clusters ($`M_i10^6\text{M}_{}`$) will not survive longer than 3.5$`\times 10^9`$ yr. This means that the disk of M51 is not the right location for young globular clusters to survive over a Hubble time.
## 8 Conclusions
We have compared the cluster population of M51 with theoretical predictions including evolutionary mass loss, cluster disruption, variable cluster formation rate and a magnitude limit. The age vs. mass diagrams of the observed cluster populations are binned to acquire two dimensional number density plots, which can be compared with simulated cluster samples. The results can be summarized as follows:
1. Artifacts introduced by our age-fitting routine do not systematically bias our sample towards young or old clusters. We present a method to correct observations of a cluster population for artifacts introduced by the age-fitting method applied.
2. The size of the largest cluster decreases with age, from 15 pc for clusters with log(Age/yr) $`<`$ 7 to 10 pc for clusters with ages around 1 Gyr. In addition, the slope of the radius distribution seems to gets steeper in time: $`\eta =2.0\pm 0.4`$ for clusters younger than log(Age/yr) = 7.5 and $`\eta =2.5\pm 0.6`$ for clusters with log(Age/yr) $`>`$ 7.5. Both these results seem to suggest that smaller clusters have a larger chance to survive. However, the radius distribution of globular clusters in our Milky Way is very similar to these values found: $`\eta =2.4\pm 0.2`$. Samples with higher spatial resolution and more clusters are needed to study the radius dependence.
3. There are more old clusters at larger distances from the galactic center. The ratio of the number of clusters in the outer parts of the galaxy (3-5 kpc) over the number of clusters in the inner part (1-3 kpc) per age bin increases with a factor of 1.8 in age (from log(Age/yr) = 6.5 to log(Age/yr) = 8.5), which is to be expected since the disruption time depends on the distance to the galactic center.
4. Assuming that the cluster disruption time depends on the initial mass of the cluster as $`t_{\mathrm{dis}}M_i^\gamma `$, and using $`\gamma =0.62`$ based on theoretical and observational studies (Lamers, Gieles & Portegies Zwart lamers05a (2005)), we find a typical disruption for a 10$`{}_{}{}^{4}\text{M}_{}^{}`$ cluster of $`t_4=1.0_{0.5}^{+0.6}\times 10^8`$ yr, where we assumed a constant cluster formation rate.
5. When $`\gamma `$ and $`t_4`$ are varied together, the value found for $`\gamma `$ is similar to that predicted by Lamers, Gieles & Portegies Zwart (lamers05a (2005)) based on observational and $`N`$-body studies. A value of $`\gamma =0.65_{0.25}^{+0.16}`$ and $`t_4=1.0_{0.35}^{+0.84}\times 10^8`$ yr are the best combination.
6. We studied the degeneracy between formation increase and disruption. Models where the cluster formation rate increases linearly in time do not affect the disruption timescale much ($`t_4`$ gets a factor 2 higher). When we include bursts at the moments of encounter with NGC 5195, the typical disruption time is also a factor of 2 higher.
7. When clusters of 10$`{}_{}{}^{4}\text{M}_{}^{}`$ are disrupted within $`2\times 10^8`$ years, and considering the power-law dependence of the disruption time scale with the initial cluster mass, even clusters with a mass of $`10^6\text{M}_{}`$ will not survive longer than 3.5 Gyr. This means that the disk of M51 is not a preferred location to form a new generation of globular clusters. This might explain why there so far are no old ($`>`$ Gyr) massive ($`>10^6\text{M}_{}`$) clusters known in the disks of spiral galaxies, although they are still forming (e.g. Westerlund 1 in the Galactic disk (Clark & Negueruela clark04 (2004)) and the young globular cluster in NGC 6946 (Larsen et al. larsen01 (2001))). |
warning/0506/math0506204.html | ar5iv | text | # Random conformal dynamical systems
## 1. Introduction
Given a probability measure on the group of homeomorphisms of a manifold, one can study the asymptotic behaviour of large composition of elements chosen randomly with respect to this measure. From the 50’s the questions of the behaviour of the random walk on the group, composition of random matrices, equidistribution of the orbits, Lyapunov exponents (if the action is by $`C^1`$ diffeomorphisms) etc. has been largely studied and understood. We are unable to review here all the history of these problems and we refer the interested reader to an excellent survey of Furman \[F\]. Nevertheless, we would like to mention here some important works in the development of this theory: Kakutani \[Kak\], Furstenberg \[Fu1\], Arnold-Krylov \[A-K\], Furstenberg-Kesten \[F-K\], Guivar’ch \[Gui\].
More recently was developed the idea that the maps could be taken from a pseudo-group, rather than a group. This has been introduced in the paper of Garnett \[Ga\] for the pseudo-group of a foliation, and then studied by Ghys \[Gh2, Gh3\], Kaimanovich \[Kai1\], Ledrappier \[Led1\] and Candel \[Can\]. Following the lines of the “Sullivan’s dictionary”, one can extend these ideas to other pseudo-groups, like for instance the one generated by an endomorphism or a correspondance.
In this work we study random compositions of the elements of a pseudo-group acting conformally on a manifold. Our results concern actions of a group by conformal transformations, conformal correspondences or transversely conformal foliations (for one dimensional dynamical systems, we suppose that the maps are of class $`C^1`$); most of them were known in the case of a group acting conformally on a compact manifold, when the conditional probabilities do not depend on the point (see \[F\]).
Observe that the existence of a measure which is invariant by every element of the pseudo-group is rather rare. In the symmetric case (i.e. when the probabilities are symmetric) we prove the following dichotomy. On a minimal subset, either is supported a probability measure which is preserved by all the elements of the pseudo-group, or the system has the property of exponential contraction: for any point and almost every random composition of elements of the pseudo-group there exists a neighborhood of the point which is contracted exponentially.
We deduce some results about the equidistribution of the orbits of the system along random compositions. In the case where the system have the property of exponential contraction (even if the system is not symmetric), we prove that the orbit of a point by almost every random compositions is distributed with respect to a unique measure. We also give examples of non symmetric systems for which the exponential contraction property and the equidistribution property are not satisfied.
### 1.1. Presentation of the results for a foliation
We begin by a survey on Garnett’s theory \[Ga\] (see also \[Can\]).
Let $``$ be a foliation of a compact manifold $`M`$, whose leaves are of class $`C^{\mathrm{}}`$, and $`g`$ a Riemannian metric on the leaves of $``$. In \[Ga\], Garnett studies the diffusion process along the leaves of $``$. Namely, the metric $`g`$ induces the Laplace-Beltrami operator along the leaves, which we denote $`\mathrm{\Delta }`$; given a continuous function $`f_0:M`$, one studies the heat equation along the leaves of the foliation
$$\frac{f}{t}=\mathrm{\Delta }f$$
with initial condition $`f(,t)=f_0`$. As it is well-known \[Cha\], because the leaves are complete and of bounded geometry, the solution to the heat equation is unique, defined for all positive time, and is expressed by convolution of the initial condition with the heat kernel $`p(x,y;t)`$. A fundamental Lemma due to Garnett (see also Candel \[Can\]), asserts that the functions $`f(,t)`$ on $`M`$ are continuous, and that the diffusion semi-group of operators $`D^t`$ defined for all $`t0`$ by
$$D^tf(0,)=f(t,)$$
acts continuously on $`C^0(M)`$.
Associated to this diffusion semi-group, Garnett considers the Brownian motion along the leaves of the foliation: this is a Markovian process with continuous time, whose trajectories stay every time in the same leaf, and whose transition probability distributions are volume forms with leafwise density given by the heat kernel. It is known that this process can be realized as a process with continuous trajectories. For any point $`x`$ in $`M`$, let $`\mathrm{\Gamma }_x`$ be the set of continuous paths parametrized by $`[0,\mathrm{})`$, starting at $`x`$, and whose image is contained in the leaf $`_x`$ passing through the point $`x`$. The space $`\mathrm{\Gamma }_x`$ is equipped with the uniform topology on compact subsets; there is a probability Borel measure induced by the Brownian motion process, expressing the probability that a trajectory occurs. This probability measure is called the Wiener measure and denoted $`W_x`$.
We recall the definition of the holonomy pseudo-group. Because the manifold $`M`$ is compact, there is a finite cover of $`M`$ by foliated box $`B_i\times T_i`$, in which the foliation $``$ is the horizontal fibration. The change of coordinates from $`B_i\times T_i`$ to $`B_j\times T_j`$ are of the form
$$(x_i,t_i)(x_j=x_j(x_i,t_i),t_j(t_i)).$$
The maps $`t_j(t_i)`$ generates a pseudo-group on the union $`T=_iT_i`$, called the holonomy pseudo-group. A measure on $`T`$ which is invariant by the holonomy pseudo-group is called a transversely invariant measure. These measures have been introduced by Schwartzman for flows \[Scm\], by Plante and Ruelle-Sullivan for foliations \[Pl, R-S\] and by Sullivan for other kind of dynamical systems \[Su1\].
Now, consider a continuous path $`\gamma `$ contained in a leaf, parametrized by a closed interval. It crosses successively the foliation boxes $`B_{i_1}\times T_{i_1},\mathrm{},B_{i_k}\times T_{i_k}`$. The composition of the associated change of transverse coordinates is by definition the holonomy map $`h_\gamma `$ corresponding to $`\gamma `$. The following result describes the asymptotic behaviour of the holonomy maps $`h_{\gamma |_{[0,t]}}`$ when $`t`$ goes to infinity, for a generic Brownian path along the leaf passing throw a point $`x`$, when the foliation is transversely conformal.
###### Theorem 1.1 (Main Theorem).
Let $``$ be a transversely conformal foliation of class $`C^1`$ of a compact manifold. Then either there exists a transversely invariant measure. Or $``$ has a finite number of minimal sets $`_1,\mathrm{},_k`$ equipped with probability measures $`\mu _1,\mathrm{},\mu _k`$, and there exists a real $`\alpha >0`$ such that:
* Contraction. For every point $`x`$ in $`M`$ and almost every leafwise Brownian path $`\gamma `$ starting at $`x`$, there is a neighborhood $`T_\gamma `$ of $`x`$ in $`T`$ and a constant $`C_\gamma >0`$, such that for every $`t>0`$, the holonomy map $`h_{\gamma |_{[0,t]}}`$ is defined on $`T_\gamma `$ and
$$|h_{\gamma |_{[0,t]}}(T_\gamma )|C_\gamma \mathrm{exp}(\alpha t).$$
* Distribution. For every point $`x`$ in $`M`$ and almost every leafwise Brownian path $`\gamma `$ starting at $`x`$, the path $`\gamma `$ tends to one of the $`_j`$ and is distributed with respect to $`\mu _j`$, in the sense that
$$\underset{t\mathrm{}}{lim}\frac{1}{t}\gamma _{}\mathrm{leb}_{[0,t]}=\mu _j,$$
where $`\mathrm{leb}_{[0,t]}`$ is the Lebesgue measure on the interval $`[0,t]`$.
* Attraction. The probability $`p_j(x)`$ that a leafwise Brownian path starting at a point $`x`$ of $`M`$ tends to $`_j`$ is a continuous leafwise harmonic function.
* Diffusion. When $`t`$ goes to infinity, the diffusions $`D^tf`$ of a continuous function $`f:M`$ converge uniformly to the function $`_jc_jp_j`$, where $`c_j=f𝑑\mu _j`$. In particular, the functions $`p_j`$ form a base in the space of continuous leafwise harmonic functions.
The existence of a transversely invariant measure for a transversely conformal foliation is a very strong condition. An ergodic component of such a measure is either supported on a compact leaf, or it is diffuse. All the examples we know of a transversely conformal foliation having a diffuse transversely invariant measure have a transverse metric which is transversely invariant. In the case of codimension one foliation of class $`C^2`$, this is an easy consequence of Sacksteder Theorem. For higher transversely conformal foliations this has been conjectured by Ghys \[Gh1\], and for codimension 3 and higher with an additional restriction of minimality this conjecture was proven by Tarquini \[Ta\].
The contraction property for group actions on the circle was studied in \[An\], \[Kaij\] and \[Kl-N\].
The distribution part of the theorem was proved by Garnett \[Ga\] in the case of the stable foliation of the geodesic flow on the unitary tangent bundle of a surface of constant negative curvature, using the contraction property and the similarity, which is straightforward in this case. This was also extended to the case of a manifold of negative variable curvature of any dimension by Ledrappier \[Led1\] (also see \[Led2\] and \[Yue\]). Finally, Hamenstädt \[Ham1\] has studied drifted Brownian motions on the stable foliation in negative curvature, and has obtained a sufficient condition for uniqueness of a harmonic measure, also giving an example (see \[Ham2\]) of non uniqueness of the harmonic measure in the drifted case.
The diffusion part of Theorem 1.1 gives examples of foliated riemannian manifolds that have non trivial continuous leafwise harmonic functions. Such an example were constructed in \[F-G\].
### 1.2. Organization of the proof.
In \[Ga\], Garnett studies the ergodic properties of the leafwise diffusion semi-group and of the leafwise Brownian motion. She introduces the notion of harmonic measure, which is a probability measure invariant by the diffusion semi-group. By the Kakutani fixed point Theorem, such a measure exists. The relation with the Brownian motion goes as follows. Consider the space $`\mathrm{\Gamma }`$ of all the continuous paths contained in a leaf of $``$. There is a semi-group $`\{\sigma _t\}_{t0}`$ of transformations of $`\mathrm{\Gamma }`$ defined for every $`t,s0`$ and every $`\gamma \mathrm{\Gamma }`$ by
$$\sigma _t(\gamma )(s)=\gamma (t+s).$$
If $`\mu `$ is a measure on $`M`$, then one consider the probability measure $`\overline{\mu }`$ on $`\mathrm{\Gamma }`$ which is defined by
$$\overline{\mu }(B)=_MW_x(B\mathrm{\Gamma }_x)𝑑\mu (x),$$
for every Borel subset $`B`$ contained in $`\mathrm{\Gamma }`$ (recall that for every $`xM`$, $`\mathrm{\Gamma }_x`$ is the set of pathes starting at $`x`$ and contained on the leaf through $`x`$, and $`W_x`$ is the Wiener measure on $`\mathrm{\Gamma }_x`$). It is straightforward to see that if $`\mu `$ is harmonic, then $`\overline{\mu }`$ is invariant by $`\{\sigma _t\}_{t0}`$, and reciprocally.
If a harmonic measure can not be written as a convex sum of different harmonic measures it is called ergodic. The Random Ergodic Theorem, due to Kakutani (\[Kak\], see also \[F\]) states that if $`\mu `$ is an ergodic harmonic measure, then $`\overline{\mu }`$ is an ergodic invariant measure of $`\{\sigma _t\}_{t0}`$. This implies in particular that for $`\mu `$-almost every point $`x`$ in $`M`$, $`W_x`$-almost every path $`\gamma \mathrm{\Gamma }_x`$ is distributed with respect to $`\mu `$. Thus, the ergodic properties of harmonic measures on the foliation $``$ can be studied via the classical ergodic theory of one dimensional semi-groups of transformations.
Most of this work deals with the Lyapunov exponent of a harmonic measure for a transversely conformal foliation $``$ of class $`C^1`$ (see \[Can, De\]). Let $`||`$ be a transverse metric. Then if $`\gamma `$ is a continuous path of $`\mathrm{\Gamma }_x`$ starting at a point $`x`$, consider the Lyapunov exponent
$$\lambda (\gamma ):=\underset{t\mathrm{}}{lim}\frac{1}{t}\mathrm{log}|Dh_{\gamma |_{[0,t]}}|,$$
when it is defined. Recall that $`h_{\gamma |_{[0,t]}}`$ is the holonomy map from a transversal $`T_{\gamma (0)}`$ passing through $`x=\gamma (0)`$ to a transversal $`T_{\gamma (t)}`$ passing through $`\gamma (t)`$. By the Birkhoff Ergodic Theorem and the Random Ergodic Theorem, if $`\mu `$ is an ergodic harmonic measure, then for $`\mu `$-almost every point $`x`$, and $`W_x`$-almost every Brownian path $`\gamma `$ starting at $`x`$, the Lyapunov exponent of $`\gamma `$ converges to a number depending only on $`\mu `$; we call it the Lyapunov exponent of the measure $`\mu `$ and denote it by $`\lambda (\mu )`$.
We begin by the study of the case where the Lyapunov exponent is negative.
###### Theorem A.
Let $``$ be a transversely conformal foliation of class $`C^1`$ of a compact manifold, and suppose that on a minimal set $``$ is supported a harmonic ergodic measure $`\mu `$ with negative Lyapunov exponent. Then the following properties are satisfied:
* Contraction. Let $`\alpha `$, $`0<\alpha <|\lambda (\mu )|`$ be chosen. Then for any $`x`$, and almost every Brownian path $`\gamma \mathrm{\Gamma }_x`$, there exist a transversal $`T_\gamma `$ at $`x`$ and a constant $`C_\gamma >0`$ such that for every $`t>0`$, the holonomy map $`h_{\gamma |_{[0,t]}}`$ is defined on $`T_\gamma `$, and
$$|h_{\gamma |_{[0,t]}}(T_\gamma )|C_\gamma \mathrm{exp}(\alpha t).$$
* Unique ergodicity. For any point $`x`$, almost every Brownian path starting at $`x`$ is distributed with respect to $`\mu `$. Thus $`\mu `$ is the unique harmonic measure on $``$.
* Diffusion. The diffusions $`D^tf`$ of a continuous function $`f:`$ converge uniformly to the constant function $`_{}f𝑑\mu `$.
* Attraction. Suppose $`M`$, and let $`p_{}(x)`$ be the probability that a Brownian path starting at $`x`$ tends to $``$, is distributed with respect to $`\mu `$, and contracts exponentially a transversal at $`x`$. Then $`p_{}`$ is lower semi-continuous and leafwise harmonic. In particular, $`p_{}`$ is bounded from below be a positive constant in some neighborhood of $``$.
Theorem A is proved in section 2. The idea is the following. A lemma of contraction (Lemma 2.2) implies, together with the fact that the Lyapunov exponent is negative, that for $`\mu `$-almost every point $`x`$, almost every Brownian path starting at $`x`$ contracts a transverse ball exponentially. We prove that if $`\overline{x}`$ is a point which is close to such a point $`x`$, then there is a similarity between the Brownian motions on the leaf $`L_{\overline{x}}`$ and on $`L_x`$. This comes from the fact that for a lot of Brownian pathes on $`L_x`$, the leaves approach each other exponentially. All the properties annouced in the Theorem A are deduced from this property of similarity.
###### Remark 1.2.
Theorem A is also true if the foliation is singular, but the minimal set does not contain any singularity. In particular, our result applies for singular holomorphic foliations on complex compact surfaces. For instance, we prove that if $``$ is a minimal subset of a holomorphic foliation of the complex projective plane, then on $``$ is supported a unique harmonic measure and the Lyapunov exponent is negative. This has been recently proved by Fornaess and Sibony for a lamination by holomorphic curves of class $`C^1`$ contained in the complex projective plane \[F-S\].
In section 3, we prove that there is a dichotomy between the case of negative Lyapunov exponent and the case where there exists a transversely invariant measure:
###### Theorem B.
Let $``$ be a transversely conformal foliation of a compact manifold. Then on a minimal subset, either there exists a transversely invariant measure, or the harmonic measure is unique and the Lyapunov exponent is negative.
In the case of a group of diffeomorphisms, the first result of this kind was proved by Furstenberg \[Fu1\]: if $`G`$ is an irreducible subgroup of projective transformations of $`P^n`$, equipped with a probability measure of finite first moment, then it has the contraction property. For a group of diffeomorphisms of class $`C^1`$ of an arbitrary compact manifold, the dichotomy “there exists a measure which is invariant by all the elements of the group or there is a stationary measure with negative sum of the Lyapunov exponents” was proved by Baxendale \[Ba\].
From Theorem B and a Theorem of Candel \[Can\], we obtain the following:
###### Corollary 1.3.
A minimal subset of a transversely conformal foliation of a compact manifold (if the codimension is one, the foliation is supposed of class $`C^1`$) carries an invariant measure, or there is a loop contained in a leaf with hyperbolic holonomy.
Our main theorem is a consequence of Theorem A and B. At the end of section 3, we prove the main theorem.
In \[Can\], Candel extends Garnett’s theory to the case of a non symmetric Laplace operator on the leaves of a foliation. For this kind of processes, Theorem A is still valid: negative Lyapunov exponent implies unique ergodicity, contraction etc. However, in the non-symmetric case Lyapunov exponent can be positive, the dynamics not uniquely ergodic, even if the foliation is minimal. Such an example is presented in Section 3.4.
The last part is a tentative to prove unique-ergodicity for a Laplace operator whose drift vector field preserves the Riemannian volume. We prove this when the foliation together with the Laplace operator are similar. A foliation equipped with a Laplace operator on the leaves is similar if there exists a transverse foliation which leaves the
Laplace operator invariant.
###### Theorem C.
Let $`(,\mathrm{\Delta })`$ be a similar codimension one foliation of a compact manifold, which is transversely continuous. Suppose that the operator $`\mathrm{\Delta }`$ is obtained by drifting the Laplacian of a Riemannian metric by a vector field that preserves the volume. Then on every minimal subset is supported a unique harmonic measure. Moreover, if $`\mathrm{\Delta }`$ is symmetric, every ergodic harmonic measure is supported on a minimal set.
The idea of the proof of Theorem C is, as in Theorem A, based on the fact that the leaves through close points stay close in a lot of “directions”. This idea is due to Thurston (\[Th\], see also \[C-D, Fe\]). Because we already know that the Brownian motion on different leaves are similar, such a property implies unique-ergodicity. This property is proved by constructing a harmonic transverse distance and to use the Martingale Theorem. This is done in section 4.
### 1.3. Other examples of dynamical systems
Following the lines of the “Sullivan’s dictionary”, our results are still valid for discrete conformal dynamical systems. The idea is that, instead of considering a foliation by smooth manifolds, one can consider foliations by graphs.
Let $`M`$ be a compact manifold together with a conformal structure of class $`C^1`$, and $`\mathrm{\Gamma }`$ be a finitely generated pseudo-group of conformal transformations of $`M`$. Then for any $`x`$ in $`M`$, we denote by $`O(x)`$ the orbit of $`x`$ under the action of the elements of $`\mathrm{\Gamma }`$. Given a symmetric system of generators of $`\mathrm{\Gamma }`$, we consider a distance on every orbit $`O(x)`$: the distance between $`x`$ and a point $`y`$ in $`O(x)`$ is the minimal number of elements of the system of generators that is necessary to map $`x`$ to $`y`$.
Consider a family $`\{\mu _x\}_{xM}`$ of probability measures on $`M`$ whose support is the orbit of $`x`$. We ask that for any element $`\gamma `$ of the pseudo-group $`G`$, the function
$$x\mathrm{dom}(\gamma )\mu _x(\gamma (x))(0,1)$$
is Hölder. The diffusion operator acts on the space of continuous functions by the formula
$$Df(x)=_{O_x}f𝑑\mu _x,$$
for any $`fC^0(M)`$ and any $`xM`$. An invariant measure by the diffusion semi-group always exists, and is called a stationary measure.
Associated to this diffusion process on the leaves, we consider the Markov process induced by the measures $`\mu _x`$ on the orbits of $`\mathrm{\Gamma }`$, which is the discrete analog of the Brownian motion on the leaves of a foliation. We define the Lyapunov exponent when the following finiteness hypothesis holds:
$$xM,\alpha >0,_{O_x}\mathrm{exp}(\alpha d(x,y))𝑑\mu _x(y)<\mathrm{}.$$
Our results extends to the context of a pseudo-group with the following analogies:
* the leaves together with the Brownian motion process is replaced by the orbit of a point together with the Markovian process with transition probabilities $`p(x,y)=_{\gamma (x)=y}\mu _x(\gamma (x))`$.
* the symmetry condition in the discrete case means that there exists a Hölder function $`v:M`$ such that for any points $`x,yM`$:
$$p(x,y)\mathrm{exp}(v(x))=p(y,x)\mathrm{exp}(v(y)).$$
The important examples are the pseudo-group generated by the action of a group on $`𝐒^n`$ by conformal transformations, or the action of an affine conformal correspondence on a torus $`𝐓^n`$.
## 2. Negative Lyapunov exponent
In this section, we are going to prove Theorem A for a codimension one foliation of class $`C^1`$. In this case, we can find a transversal foliation of dimension one and of class $`C^1`$, and this will simplify the proof. Such a foliation does not necessarily exist in the case of higher codimension, but at the appendix (Section 5.3) we explain how to adapt the proof. Finally, we may suppose $``$ to be transversely orientable (and we do so from this moment), as it is always true up to a 2-folded cover, and passing a finite cover does not change our results.
Let $``$ be a codimension one foliation of class $`C^1`$ of a compact manifold $`M`$, and $`𝒢`$ a transverse foliation of class $`C^1`$. Consider some minimal subset $`M`$. Let us suppose that for some ergodic harmonic measure $`\mu `$ supported on $``$, we have $`\lambda (\mu )<0`$.
### 2.1. Contraction
The goal of this paragraph is to prove that for generic Brownian paths, the holonomy contracts a transverse interval exponentially, which is true infinitesimally:
###### Proposition 2.1.
Suppose that there exists an ergodic harmonic measure $`\mu `$ on $`M`$, such that $`\lambda (\mu )<0`$. Let $`\alpha >0,\alpha <|\lambda (\mu )|`$. Then for $`\mu `$–almost every $`xM`$ and for $`W_x`$-almost every $`\gamma \mathrm{\Gamma }_x`$, there exists a transversal neighborhood $`I`$ and a constant $`C>0`$, such that
$$t>0\left|h_{\gamma |_{[0,t]}}(I)\right|<C|I|e^{\alpha t}.$$
In particular, all the holonomy maps $`h_{\gamma |_{[0,t]}}`$ are defined in the same transversal neighborhood $`I`$ and, as $`t\mathrm{}`$, this neighborhood is exponentially contracted.
To this end, we have to connect the derivatives of holonomy maps in one point $`x`$ and the diameter of the images of the transversal neighborhood $`I`$. The following Contraction Lemma is in the “folklore”. For the completeness of the text, we present here both its statement and proof. This lemma extends the Distortion Lemmas, used in $`C^2`$ case by Schwartz \[Sch\], Denjoy \[Den\] and Sacksteder \[Sa\], in $`C^{1+\tau }`$ by Sullivan \[Su2\] and Hurder \[Hur\].
###### Lemma 2.2 (Contraction Lemma).
Let $`x_0,x_1,\mathrm{}`$ be points in $``$, $`I_j,I_j=U_\epsilon (x_j)`$ be their $`\epsilon `$-neighborhoods. Let
$$h_j:I_j,h_j(x_j)=x_{j+1},j=0,1,2,$$
be $`C^1`$-diffeomorphisms onto their image. Let $`\overline{f}_j(y)=h_j(y+x_j)x_{j+1}`$ be diffeomorphisms of $`U_\epsilon (0)`$, and suppose, that these diffeomorphisms are bounded in the $`C^1`$ topology.
Denote
$$F_n=h_n\mathrm{}h_1:I_1,n𝐍,$$
and suppose that
(2.1)
$$\underset{n\mathrm{}}{lim\; sup}\frac{1}{n}\mathrm{log}F_n^{}(x_0)=\lambda <0,$$
and let $`\alpha >0,\alpha <|\lambda |`$. Then there exist such an $`\epsilon _1>0`$ and a constant $`C`$, that for any interval $`JI_0`$, such that $`x_0J,|J|<\epsilon _1`$ all the compositions $`F_n`$ are defined on $`J`$ and we have a bound
$$n|F_n(J)||J|Ce^{\alpha n}.$$
Proof. Let us choose $`\beta `$, $`\alpha <\beta <|\lambda |`$. Then, the condition (2.1) implies that the supremum
$$C=\underset{n}{sup}e^{\beta n}F_n^{}(x_0)$$
is finite. Then, for every $`n`$ we have $`F_n^{}(x_0)Ce^{\beta n}`$.
Note, that due to pre-compactness property the logatithms of the derivatives $`\mathrm{log}h_n^{}`$ are uniformly continuous. Thus, there exists $`\epsilon _0>0`$, such that for every $`n`$ and for every $`y,zI_n`$, $`|yz|<\epsilon _0`$ we have $`h_n^{}(y)/h_n^{}(z)<e^{\beta \alpha }`$.
Now, let $`\epsilon _1=\mathrm{min}(\epsilon _0,\epsilon )/C`$, and let $`J`$ be an interval of length less than $`\epsilon `$, containing $`x_0`$. We are going to prove the following statement: for every $`n`$, the length of $`|F_n(J)|<Ce^{\alpha n}|J|`$. In fact, by the mean value theorem $`|F_n(J)|=F_n^{}(y)|J|`$ for some point $`yJ`$. Now,
(2.2)
$$\frac{F_n^{}(y)}{F_n^{}(x)}=\underset{j=1}{\overset{n}{}}\frac{h_j^{}(y_{j1})}{h_j^{}(x_{j1})}<(e^{\beta \alpha })^n=e^{(\beta \alpha )n}.$$
Here $`y_j=F_j(y)`$, and the inequality $`\frac{h_j^{}(y_{j1})}{h_j^{}(x_{j1})}<e^{\beta \alpha }`$ is satisfied due to the recurrence hypothesis and the choice of $`\epsilon _0`$:
$$|x_{j1}y_{j1}||F_{j1}(J)|Ce^{\alpha n}|J|C\epsilon _1<\epsilon _0.$$
Now, from (2.2) we have:
$$F_n^{}(y)e^{(\beta \alpha )n}F_n^{}(x)e^{(\beta \alpha )n}Ce^{\beta n}=Ce^{\alpha n}.$$
Hence,
$$|F_n(J)|=F_n^{}(y)|J|<Ce^{\alpha n}|J|.$$
This proves the recurrence step, and thus the entire lemma. $`\mathrm{}`$
Now we are going to apply the Contraction Lemma to the Brownian motion on the leaves. In order to do that, we would like to decompose the holonomy map in time $`t`$ as a composition of some number $`n`$ of maps, forming a $`\mathrm{Diff}^1`$-pre-compact set, with the quotient $`n/t`$ being bounded from above and from below.
Suppose $`\delta >0`$ be given. Then to any trajectory $`\gamma \mathrm{\Gamma }_x`$, we associate a sequence of points $`(x_n)=(\stackrel{~}{\gamma }(n\delta ))`$ on the universal cover $`\stackrel{~}{}_x`$, where the path $`\stackrel{~}{\gamma }`$ is the covering path for the path $`\gamma `$, and a sequence of numbers $`k_n=[\mathrm{dist}(x_{n1},x_n)]+1`$ (here $`[z]`$ denotes the integer part of $`z`$). Let us divide a segment of shortest geodesic line, joining $`x_{n1}`$ and $`x_n`$, in $`k_n`$ equal parts; let us denote the vertices of this partition by $`y_n^0=x_{n1},y_n^1,\mathrm{},y_n^{k_n1},y_n^{k_n}=x_n`$. Then, the holonomy map $`h_{x_0x_n}`$ between $`x_0`$ and $`x_n`$ can be written as
$$h_{x_0x_n}=h_{x_{n1}x_n}h_{x_{n2}x_{n1}}\mathrm{}h_{x_0x_1}$$
and thus as
(2.3)
$$h_{x_0x_n}=(h_{y_n^{k_n1}y_n^{k_n}}\mathrm{}h_{y_n^0y_n^1})\mathrm{}(h_{y_1^{k_n1}y_1^{k_n}}\mathrm{}h_{y_1^0y_1^1}).$$
The total number of the maps in the right hand side of (2.3) is $`K_n=k_1+\mathrm{}+k_n`$. The following lemma is some form of the Large Numbers Law. Though it is rather clear that such statement should take place, its rigorous proof is rather long, and we have put it in Section 5.
###### Lemma 2.3.
There exists a constant $`c>0`$, such that for any $`xM`$ for $`W_x`$–almost every path $`\gamma \mathrm{\Gamma }_x`$ we have $`K_n/n<c`$ for all $`n`$ sufficiently big.
Thus, the discretization of the Brownian motion is “quasi-preserving” the time: the number of terms in the right-hand side of the representation (2.3) is comparable with the time passed.
Proof of Proposition 2.1. Suppose that a point $`xM`$ is such that for almost every path $`\gamma \mathrm{\Gamma }_x`$ we have
(2.4)
$$\underset{t\mathrm{}}{lim}\frac{1}{t}\mathrm{log}h_{\gamma |_{[0,t]}}^{}(x)=\lambda (\mu )<0.$$
Let us show that for this point the conclusion of Proposition 2.1 holds. This will prove the Proposition — as $`\mu `$–almost every point $`x`$ satisfies (2.4).
Recall that $`\alpha <|\lambda |`$. As it follows from (2.4), for $`W_x`$–almost every path $`\gamma \mathrm{\Gamma }_x`$ there exists a constant $`C_0>0`$, such that for every $`t>0`$
(2.5)
$$h_{\gamma |_{[0,t]}}^{}(x)<C_0e^{\alpha t}.$$
Let us consider a discretization $`(x_n,k_n)`$ of such path $`\gamma `$. As it follows from Proposition 2.3, for almost every path $`\gamma `$ we have also
(2.6)
$$N:n>N\frac{K_n}{n}<c.$$
Suppose that for $`\gamma `$ both (2.5) and (2.6) are satisfied. Choosing some $`c^{}>c`$, we may suppose that for all $`n`$ we have $`K_n<c^{}n`$. As it follows from (2.3), for every $`n`$ the holonomy map $`h_{x_0x_n}`$ can be written as a composition of $`K_n<c^{}n`$ maps $`h_{y_l^jy_l^{j+1}}`$, each one being a holonomy between two points at the distance at most $`1`$. The set of holonomy maps along paths of length at most 1 is pre-compact (it is a continuous image of a compact set), and the derivative of such composition at the point $`x`$ is less then $`C_0e^{\alpha n}<C_0e^{\alpha K_n/c^{}}`$, thus, we still have an exponential decrease of derivatives at $`x`$ with respect to the number of maps. The application of the Lemma 2.2 concludes the proof. $`\mathrm{}`$
### 2.2. Similarity of the Brownian motions on different leaves.
Suppose that the point $`x`$ is typical in the sense of Proposition 2.1, and that almost every Brownian path starting at $`x`$ is distributed with respect to $`\mu `$. Then there exist a transversal interval $`I`$, and constants $`C_0`$, $`\alpha >0`$, such that the set
$$E_x=\{\gamma \mathrm{\Gamma }_xt\left|h_{\gamma |_{[0,t]}}(I)\right|<C_0\mathrm{exp}(\alpha t)|I|\}$$
has positive Wiener measure: $`W_x(E_x)>0`$ (in fact, this probability can be made arbitrary close to $`1`$ by choosing sufficiently small $`I`$ and sufficiently big $`C_0`$). Let us fix such $`I`$ and $`C_0`$.
If $`\overline{x}`$ is a point close to $`x`$, consider the set $`E_{\overline{x}}`$ of Brownian paths $`\overline{\gamma }\mathrm{\Gamma }_{\overline{x}}`$ which are exponentially asymptotic to a path $`\gamma `$ of $`\mathrm{\Gamma }_{\overline{x}}`$:
$$t0,d(\gamma (t),\overline{\gamma }(t))<C_0\mathrm{exp}(\alpha t)|I|.$$
For instance, a path of $`E_{\overline{x}}`$ can be constructed from a path of $`E_x`$ by following the foliation $`𝒢`$.
###### Lemma 2.4 (Similarity of the Brownian motions).
There exists a neighborhood $`U`$ of $`x`$ such that for any $`\overline{x}U`$, the Wiener measure of $`E_{\overline{x}}`$ is positive. Moreover, it can be made arbitrarily close to $`1`$ by choosing sufficiently small $`I`$, large $`C_0`$ and small $`U`$. Finally, $`W_{\overline{x}}`$-almost every path $`\gamma E_{\overline{x}}`$ is distributed with respect to $`\mu `$.
We are going to outline the proof of this Lemma. However, a formal realization of the ideas faces some difficulties and becomes very technical. Thus, we have postponed it until Appendix (section 5).
First, we observe that the lemma is simple to prove when the foliation is similar, meaning that the foliation $`𝒢`$ preserves the Laplace operator. In this case, the $`𝒢`$-along holonomy preserves the metric and thus translates the Brownian motion on an initial leaf to the Brownian motion of the range leaf. In particular, the probabilities $`W_x(E_x)`$ and $`W_{\overline{x}}(E_{\overline{x}})`$ coincide if the points $`x`$ and $`\overline{x}`$ are in the same $`𝒢`$-leaf. On the other hand, if we consider Brownian motions starting at a point $`y`$ and at some point $`\overline{y}_y`$ close enough to $`y`$, the distribution of their values at the moment $`\delta `$ will be absolutely continuous with respect to each other, and the density will be close to 1. Thus, the probabilities of any tail-type properties (in particular, of belonging to $`E_y`$ and $`E_{\overline{y}}`$, being distributed with respect to $`\mu `$, etc.) are close enough. The $`W_x(E_x)>0`$ property implies the same for the points close enough on the same leaf of $``$ and for the points close enough on the same leaf of $`𝒢`$. Thus, this property is satisfied in some neighborhood of $`x`$ (as $``$ and $`𝒢`$ are transversal). This proves the lemma for similar foliations.
In the general case, the $``$-leafwise implication is still valid. Unfortunately, it is much more difficult to prove the $`𝒢`$-along implication, when the Riemannian metric is not invariant by the foliation $`𝒢`$. Let us suppose that $`\overline{x}𝒢_x`$. Denote by $`\mathrm{\Phi }_{x,\overline{x}}:E_xE_{\overline{x}}`$ the holonomy map along the transversal foliation $`𝒢`$. Except for similar foliations, $`\mathrm{\Phi }_{x,\overline{x}}`$ does not translate the Wiener measure $`W_x`$ on $`E_x`$ to a measure which is absolutely continuous with respect to the Wiener measure $`W_{\overline{x}}`$ on $`E_{\overline{x}}`$. This effect comes from small movements — typical path of the Brownian motion, being considered on arbitrary small interval of time, allows to reconstruct the Riemannian metrics (on its support). Thus, we are going to pass from the Brownian paths to their discretization — as it was already done in the previous paragraph.
Let $`\delta >0`$ be given. Let us define the discretization map $`F^\delta :E_x(\stackrel{~}{}_x)^{\mathrm{}}`$ as
(2.7)
$$F(\gamma )=\{\stackrel{~}{\gamma }(n\delta )\}_{n=0}^{\mathrm{}}.$$
Denote by $`E_x^\delta `$ the image of $`E_x`$ under $`F^\delta `$, and by $`W_x^\delta `$ measure on $`E_x^\delta `$ that is the image of the Wiener measure on $`\mathrm{\Gamma }_x`$, restricted on $`E_x`$, under $`F^\delta `$. We claim that if $`\overline{x}𝒢_x`$ is sufficiently close to $`x`$, then the induced map
$$\mathrm{\Phi }_{x,\overline{x}}^\delta :(E_x^\delta ,W_x^\delta )(E_{\overline{x}}^\delta ,W_{\overline{x}}^\delta )$$
is absolutely continuous and its Radon-Nykodym derivative is uniformly close to $`1`$ on a set of a large measure. This comes from the fact that for a finite number of steps of discretization the density can be written explicitly: the density is the product of quotients of heat kernels on these two leaves. The step of discretization is constant and equals $`\delta `$, and the leaves approach each other along the trajectories of $`E_x`$ exponentially. Thus, it is natural to expect that such product would converge. Moreover, the closer initial points $`x`$ and $`\overline{x}`$ are, the closer to $`1`$ will be the product. These facts imply that the probabilities $`W_x(E_x)`$ and $`W_{\overline{x}}(E_{\overline{x}})`$ are sufficiently close and thus imply the Lemma; their rigorous proof can be found in the appendix (Section 5).
To simplify the notations, we note $`\sigma =\sigma _\delta :\mathrm{\Gamma }\mathrm{\Gamma }`$:
$$t0\sigma (\gamma )(t)=\gamma (t+\delta ).$$
###### Corollary 2.5 (Similarity of Brownian motions).
For any $`y`$ for $`W_y`$-almost every path $`\gamma \mathrm{\Gamma }_y`$ there exists $`n`$, such that $`\gamma (n\delta )U`$, $`\sigma ^n(\gamma )E_{\gamma (n\delta )}`$, and $`\gamma `$ is distributed with respect to $`\mu `$.
Proof. Let us consider the map $`D:\mathrm{\Gamma }\mathrm{\Gamma }`$, defined in the following “algorithmic” way:
$$D=D_1D_2,$$
where $`D_2(\gamma )=\sigma ^n(\gamma )`$, $`n=inf\{j\gamma (j\delta )U\}`$,
$$D_1(\gamma )=\{\begin{array}{cc}\text{STOP},\hfill & \text{if }\gamma E_{\gamma (0)},\hfill \\ \sigma ^n(\gamma ),\hfill & n=inf\{j|h_{\gamma |_{[0,j\delta ]}}(I)|Ce^{\alpha j\delta }|I|\}.\hfill \end{array}$$
Note, that $`D_2(\gamma )`$ is defined on $`W_z`$-almost every $`\gamma \mathrm{\Gamma }_z`$ for every $`z`$ due to the minimality of $``$ and due to the Markovian property of the Brownian motion: on every step of discretization, the probability of hitting $`U`$ is positive and bounded from below; thus the probability of never hitting $`U`$ is 0.
Now, note that for every $`z`$ the probability of $`D`$ returning “STOP” on $`\gamma \mathrm{\Gamma }_z`$ is bounded from below due to Lemma 2.4. Also, due to the Markovian property for every $`z,z^{}`$ the conditional distribution of $`D(\gamma )`$, $`\gamma \mathrm{\Gamma }_z`$ with respect to the condition $`D(\gamma )(0)=z^{}`$, coincides with $`W_z^{}`$. Thus, the probability of $`D`$ not stopping on $`\gamma `$ in $`k`$ iterations tends to zero exponentially. In particular, $`W_z`$-almost surely some iteration $`D^k(\gamma )`$ stops, which means, that $`D_2(D^{k1}(\gamma )E_{\overline{z}}`$, where $`\overline{z}=D_2(D^{k1}(\gamma ))(0)`$. Together with the definitions of $`D_1`$ and $`D_2`$, this proves the first part of the corollary.
Finally, recall that for any $`n`$ the conditional distribution of $`\sigma ^n(\gamma )`$ with respect to the condition $`\gamma (n\delta )=z`$ coincides with $`W_z`$. Lemma 2.4 states, that for every $`zU`$ almost every path $`\gamma E_z`$ is distributed with respect to $`\mu `$, hence, the constructed path $`\sigma ^n(\gamma )`$ almost surely is distributed with respect to $`\mu `$. A finite shift can not change the asymptotic distribution of a path, and so the corollary is proven. $`\mathrm{}`$
Before going into the applications of the preceeding results, we would like to make a disgression that may clarify the meaning of “similarity of the Brownian motion on the leaves”. Recall that every Riemannian manifold $`N`$ of bounded geometry has a boundary associated to the Brownian motion process on it: this is the Poisson boundary $`P(N)`$. We recall the following facts, that characterize the Poisson boundary (see \[Kai1\]). For every $`x`$ of $`N`$ there is a canonical projection $`\pi _x:\mathrm{\Gamma }_xP(N)`$. The family of probability measures $`\nu _x=(\pi _x)_{}W_x`$ depends harmonically on $`N`$, in the sense that for any bounded measurable function $`f`$ on $`P(N)`$, the function
$$P(f)(x)=_{P(N)}f𝑑\nu _x$$
is a bounded harmonic function on $`N`$. Moreover, every bounded harmonic function on $`N`$ is obtained in this way.
There is a following question, which (even if the answer is negative), in our opinion, clarifies the proof of Lemma 2.4, giving the good general idea.
###### Question 2.6.
Given two leaves $`_x`$ and $`_{\overline{x}}`$, does the argument of the proof of Lemma 2.4 allow to identify large parts of the Poisson boundaries of their universal covers $`\stackrel{~}{}_x`$ and $`\stackrel{~}{}_{\overline{x}}`$, corresponding to the couple of “directions” in which the leaves are converging exponentially to each other?
In the case of a foliation by hyperbolic surfaces of a $`3`$-manifold, Thurston has constructed the “circle at infinity” (see \[C-D, Fe\]). A leaf of the universal cover of such a foliation is isometric to the hyperbolic plane, and its boundary (as a Gromov hyperbolic space) is a topological circle. Thurston has proved that there is a natural topological identification of the circle at infinity of the leaves of the foliation on the universal cover.
###### Question 2.7.
The boundary (as a Gromov hyperbolic space) of the hyperbolic plane is also the Poisson boundary. Is it true that the topological identifications of the boundaries of Thusrton’s theorem preserve the structure of the Poisson boundary as well?
### 2.3. Proof of Theorem A
#### 2.3.1. Contraction property
Let $`x`$. Then by Corollary 2.5, for almost every Brownian path $`\gamma `$ starting at $`x`$, there exists $`n`$ such that $`\gamma (n\delta )`$ belongs to $`U`$, and $`\sigma ^n\gamma E_{\gamma (n\delta )}`$. Thus starting from the time $`n\delta `$, the path $`\gamma `$ contracts a transverse interval exponentially, with exponent $`\alpha `$. The contraction property is proved.
#### 2.3.2. Unique ergodicity
Let $`x`$. Recall first that almost every Brownian path starting at $`x`$ is distributed with respect to $`\mu `$, as it is claimed by Corollary 2.5.
Let us prove that $`\mu `$ is the only harmonic measure supported on $``$. Let $`\mu ^{}`$ be an ergodic harmonic measure. Then for $`\mu ^{}`$-almost every point $`y`$, $`W_y`$-almost every Brownian path $`\gamma W_y`$ is distributed with respect to $`\mu ^{}`$. Fix one of these points $`y`$. By the preceeding argument, $`W_y`$-almost every Brownian path $`\gamma `$ is also distributed with respect to $`\mu `$. Thus $`\mu ^{}=\mu `$ and $`\mu `$ is the unique harmonic measure supported on $``$.
#### 2.3.3. Attraction
The lower semicontinuity of the function $`p_{}`$ is immediately implied by the proof of Lemma 2.4. Now, $`p_{}|_{}=1`$ due to Corollary 2.5. Thus, $`p_{}`$ is bounded from below by a positive constant in some neighborhood of $``$.
Now, due to the Markovian property of the Brownian motion, for any initial point $`x`$ the process $`p_{}(\gamma (t))`$ is a martingale. Due to Ito formula, $`p_{}`$ is leafwise harmonic.
The Theorem A is proven unless for the diffusion part, which is absolutely analogous to the proof of the diffusion part of the Main Theorem, see paragraph 3.2.
###### Remark 2.8.
In fact, for any point $`yM`$ almost all the paths $`\gamma \mathrm{\Gamma }_y`$, tending to $``$ (if such paths exist), are distributed with respect to $``$. In particular,
$$p_{}(x)=W_x(\{\gamma \mathrm{\Gamma }_x\gamma \})$$
Let us prove this. For a trajectory starting at sufficiently small distance from $``$, with the probability close to 1 this trajectory will be distributed with respect to $`\mu `$. Decompose trajectories starting at $`y`$ and tending to $``$ in two parts: finite part before arriving close to $``$ and infinite afterwards. The Markovian property implies, that the probability of the trajectory being distributed with respect to $`\mu `$ is arbitrary close to total probability of tending to $``$. Thus, as the distance of decomposition can be chosen arbitrary small, almost every trajectory, tending to $``$, is distributed with respect to $`\mu `$.
Denote by $`Attr()`$ the basin of attraction of $``$: this is the union of the leaves whose closure contains $``$. Now, it is rather clear that the function $`p_{}`$ is positive exactly in the points of $`Attr()`$.
###### Corollary 2.9.
Any harmonic ergodic measure different from $`\mu `$ has a support disjoint from $`Attr()`$.
Proof. A trajectory starting at $`yAttr()`$ tends to $``$ and is distributed with respect to $`\mu `$ with positive probability. If there exists another harmonic ergodic measure $`\mu ^{}`$ with the support not disjoint from $`Attr()`$, for $`\mu ^{}`$-almost every point $`zAttr()`$ almost all trajectories, starting at $`z`$, would be distributed with respect to $`\mu ^{}`$, not with respect to $`\mu `$. But there are no such points, which gives us the desired contradiction. $`\mathrm{}`$
### 2.4. Examples: holomorphic foliations on complex surfaces
Let $``$ be a singular holomorphic foliation of a compact complex surface $`S`$, and let $`g`$ be a hermitian metric on $`T`$. We note $`\mathrm{\Delta }`$ the Laplacian of $`g`$ along the leaves of $``$. With the use of harmonic measures, one can extend certain notions that we have for compact holomorphic curves; for instance the Euler characteristic. The following definition is due to Candel (see \[Gh4\]).
###### Definition 2.10.
Let $`ES`$ be a holomorphic line bundle over $`S`$, and $`\mu `$ a harmonic measure. The Chern-Candel class of $`E`$ against $`\mu `$ is
$$c_1(E,\mu ):=\frac{1}{2\pi }_S\mathrm{curvature}(||)d\mu ,$$
where $`||`$ is a hermitian metric on $`E`$ of class $`C^2`$. Recall that the curvature of a hermitian metric is
$$\mathrm{curvature}(||):=2\mathrm{\Delta }\mathrm{log}|s|,$$
where $`s`$ is a local non vanishing holomorphic section of $`E`$. Because the curvature of two different smooth hermitian metrics on $`E`$ differs by the Laplacian along $``$ of a smooth function, the Chern-Candel class of $`E`$ does not depend on the choice of the hermitian metric. The Euler characteristic of $`T`$ is the first Chern class of the tangent bundle of $``$.
The following lemma expresses the Lyapunov exponent in algebraic terms.
###### Lemma 2.11.
Let $``$ be a closed minimal subset, which does not contain singularities of $``$. Let $`g`$ be a hermitian metric on $`T`$ and $`\mu `$ a harmonic measure supported on $``$. Then
$$\lambda (\mu )=\pi c_1(N_{},\mu ).$$
Proof. Let $`||`$ be a smooth conformal transverse metric of the foliation. Let us compute the curvature of the normal bundle $`N_{}`$. Let $`(z,t)`$ some local coordinates where the foliation is defined by $`dt=0`$. The section $`\frac{}{t}`$ induces a non vanishing holomorphic section of the normal bundle. Thus, the curvature of $`N_{}`$ is
$$\mathrm{curvature}(||)=2\mathrm{\Delta }\mathrm{log}|\frac{}{t}|.$$
One gets
$$c_1(N_{},T_{g,\mu })=\frac{1}{\pi }\mathrm{\Delta }\mathrm{log}|\frac{}{t}|d\mu ,$$
and the lemma is proved by applying Lemma 3.1. $`\mathrm{}`$
###### Theorem 2.12.
Let $``$ be a singular holomorphic foliation of a complex surface, and $``$ be an exceptional minimal subset of $``$. Suppose that the normal bundle of $``$ has a metric of positive curvature along $``$. Then on $``$ is supported a unique harmonic measure of negative Lyapunov exponent.
Proof. Let $`\mu `$ be an ergodic harmonic measure supported on $``$. The normal bundle of $``$ has a metric of positive curvature. Thus, by Lemma 2.11, the Lyapunov exponent of $`\mu `$ is negative. The theorem is a corollary of Theorem A. $`\mathrm{}`$
Recently Fornaess-Sibony proved that on a compact lamination by holomorphic curves of class $`C^1`$ of $`P^2`$, there is a unique harmonic measure \[F-S\]. We recover this result here for an exceptional minimal set of a singular holomorphic foliation of $`P^2`$.
###### Corollary 2.13.
An exceptional minimal subset of a singular holomorphic foliation of $`P^2`$ carries a unique harmonic measure which is of negative Lyapunov exponent.
Proof. The normal bundle of a singular holomorphic foliation of $`P^2`$ is $`O(d+2)`$ where $`d>0`$ is the degree. This bundle has a metric of positive curvature everywhere. $`\mathrm{}`$
## 3. Symmetric case
A Laplace operator is called symmetric if it is the Laplace-Beltrami operator of a Riemannian metric. In this section we consider a foliation equipped with a Riemannian metric on the leaves, and the Laplace-Beltrami operator $`\mathrm{\Delta }`$ along the leaves.
### 3.1. The dichotomy: proof of Theorem B
In this paragraph we prove Theorem B: on a minimal subset of a transversely conformal foliation is supported a unique harmonic measure with negative Lyapunov exponent, or there is a transversely invariant measure.
If there exists a harmonic measure for which the Lyapunov exponent is negative, then Theorem A shows that this is the unique harmonic measure. Thus we will prove that if the Lyapunov of any harmonic measure is non negative, then there exists a transversely invariant measure.
We use an integral formula which expresses the Lyapunov exponent, and which has been founded in \[Can, De\]. Let $`||`$ be a transverse conformal metric. In a foliation box we consider a transverse vector field $`u`$ which is invariant by the holonomy, and define $`\phi :=|u|`$. Note that $`\phi `$ is well-defined up to the multiplication by a leafwise constant function. Thus the function $`f=\mathrm{\Delta }\mathrm{log}\phi `$ is a well-defined continuous function on $`M`$.
###### Lemma 3.1.
For any ergodic harmonic measure $`\mu `$ on $`M`$, $`\lambda (\mu )=_Mf𝑑\mu `$.
Proof. For every $`t0`$, let $`L_t:\mathrm{\Gamma }`$ be the functional defined by:
$$L_t(\gamma )=\mathrm{log}|Dh_{\gamma |_{[0,t]}}|.$$
The family $`\{L_t\}_{t0}`$ is a cocycle with respect to the shift semi-group, in the sense that for any $`s,t0`$:
$$L_{t+s}=L_t+L_s\sigma _t.$$
Thus the integrals $`_\mathrm{\Gamma }L_t𝑑\overline{\mu }`$ depends linearly of $`t`$. Because by definition the limit
$$\lambda (\gamma )=\underset{t\mathrm{}}{lim}\frac{L_t(\gamma )}{t}$$
exists for $`\overline{\mu }`$-almost every path and is equal to the Lyapunov exponent of $`\mu `$, we have for every $`t0`$:
$$_\mathrm{\Gamma }L_t𝑑\overline{\mu }=t\lambda (\mu ).$$
Now, let $`x,y`$ be two points in the universal cover $`\stackrel{~}{L}`$ of a leaf $`L`$. Define the cocycle
$$c(x,y):=\mathrm{log}h_{x,y}^{},$$
where $`h_{x,y}`$ is the holonomy between $`x`$ and $`y`$. Let
$$\lambda _t(x):=_{\stackrel{~}{L}}p(x,y;t)c(x,y)\mathrm{vol}_g(y).$$
The work of Garnett (\[Ga, p. 288, Fact 1\]) shows that $`\lambda _t`$ is a continuous function on $`M`$. Moreover we have the formula
(3.1)
$$t\lambda (\mu )=_\mathrm{\Gamma }L_t𝑑\overline{\mu }=_M\lambda _t(x)𝑑\mu (x).$$
Observe that $`\phi `$ is a function defined up to multiplication by a constant on $`\stackrel{~}{L}`$ and that by definition $`\phi (y)/\phi (x)=h_{x,y}^{}`$. Thus we get $`c(x,y)=\mathrm{log}\phi (y)\mathrm{log}\phi (x)`$ and taking the derivative at $`t=0`$:
(3.2)
$$\frac{d\lambda _t(x)}{dt}|_{t=0}=\mathrm{\Delta }\mathrm{log}\phi (x).$$
Differentiating (3.1) at $`t=0`$ and substituting (3.2), we obtain the desired result. $`\mathrm{}`$
We will also use another formula which expresses the conservation of mass of the diffusion semi-group. Let $`v`$ be the volume form on $`M`$, induced by the leafwise volume form $`\mathrm{vol}_g`$ and the transverse metric.
###### Lemma 3.2.
$`_M(\mathrm{\Delta }\phi +|\phi |^2)v=0`$.
Proof. Let $`v_t`$ be the transverse volume form induced by the transverse metric. Then, in a foliation box $`B\times T`$, there is a volume form $`\theta `$ on $`T`$ such that
$$v_t=\mathrm{exp}(q\phi )\theta ,$$
where $`q`$ is the codimension of $``$. This is by definition of $`\phi `$.
Consider a partition of unity: $`1=_if_i`$, where the support of each function $`f_i`$ is contained in a foliation box $`B_i\times T_i`$. By Green’s formula:
$$\begin{array}{c}_M\mathrm{\Delta }f_iv=_{B_i\times T_i}(\mathrm{\Delta }f_i)\mathrm{exp}(\phi _i)\mathrm{vol}_g\theta _i=\hfill \\ \hfill =_{B_i\times T_i}f_i(\mathrm{\Delta }\mathrm{exp}(\phi _i))\mathrm{vol}_g\theta _i=_{B_i\times T_i}f_i(\mathrm{\Delta }\phi +|\phi |^2)\mathrm{exp}(\phi )\mathrm{vol}_g\theta _i=\\ \hfill =_Mf_i(\mathrm{\Delta }\phi +|\phi |^2)v.\end{array}$$
By summing these equalities and using the fact that $`\mathrm{\Delta }1=0`$, we obtain
$$\begin{array}{c}0=_M\mathrm{\Delta }1v=\underset{i}{}_M\mathrm{\Delta }f_iv=\hfill \\ \hfill =\underset{i}{}_Mf_i(\mathrm{\Delta }\phi +|\phi |^2)v=_M(\mathrm{\Delta }\phi +|\phi |^2)v.\end{array}$$
Lemma 3.2 is proved. $`\mathrm{}`$
#### 3.1.1. Proof of Theorem B when $`=M`$
By the integral definition of a harmonic measure (see \[Ga, Lemma B, p. 294\]) and Lemma 3.1, the fact that the Lyapunov exponent of every harmonic measure is non negative is trivially satisfied if the function $`f`$ is bounded from below by the laplacian along the leaves of a smooth function $`h`$:
$$f\mathrm{\Delta }h.$$
Let us study this case as a relevant example. Let $`v_t`$ be the transverse volume form induced by the transverse metric $`||`$, and $`v_t^{}`$ be defined by $`v_t^{}=\mathrm{exp}(qh)v_t`$, where $`q`$ is the codimension of $``$. We claim that $`v_t^{}`$ is transversely invariant, or, what is the same, the measure $`\mu =\mathrm{vol}_gv_t^{}`$ is totally invariant. By Lemma 3.2, we have
$$_M(\mathrm{\Delta }\phi ^{}+|\phi ^{}|^2)𝑑\mu =0.$$
Because $`\phi ^{}`$ is subharmonic, this implies that $`\phi ^{}`$ is locally constant along the leaves. Thus $`\mu `$ is a totally invariant measure.
###### Remark 3.3.
A more geometric proof of this goes as follows: we consider the leafwise gradient of the function $`\mathrm{log}\phi ^{}`$, which is well-defined everywhere. Due to sub-harmonicity of $`\phi ^{}`$ it dilates leafwise volume. By definition it also dilates the transverse volume. Thus, the total volume must increase everywhere. The only way is that the function $`\phi ^{}`$ must be leafwise constant, meaning that the measure $`\mu `$ is totally invariant.
Unfortunately, there exists a diffeomorphism of the circle, which is minimal, and whose invariant transverse measure is singular \[K-H\]. Therefore, we can not hope to solve the functional inequality
$$f\mathrm{\Delta }_{}h$$
in general. However, it is still possible to solve this inequality “approximatively”. The following result is due to Ghys \[Gh2, Gh3\], and the proof is based on the use of the Hahn-Banach theorem; this idea goes back to the famous paper \[Su1\] by Sullivan on the foliation cycles.
###### Lemma 3.4.
Let $`f:M`$ be a continuous function such that
$$_Mf𝑑\mu 0$$
for every harmonic measure $`\mu `$. Then there exists a sequence of smooth functions $`\psi _n:M`$ such that uniformly,
$$\underset{n\mathrm{}}{lim\; inf}f\mathrm{\Delta }_{}\psi _n0.$$
Proof. In the Banach space $`C^0(M)`$, consider the closed subspace $`E`$ of uniform limit of leafwise Laplacian of smooth functions, the cone $`𝒞`$ of everywhere positive functions. Let $`F=C^0(M)/E`$ and $`\overline{𝒞}F`$ be the closure of the image of the cone $`𝒞`$ under the natural projection $`C^0(M)F`$. The conclusion of Lemma 3.4 is equivalent to the fact that the image of $`f`$ in $`F`$ is in $`\overline{𝒞}`$. Suppose it is not the case. Then, by Hahn-Banach separation theorem, there exists a linear functional which is non-negative on $`\overline{𝒞}`$ and negative on the image of $`f`$. Such a linear functional is by definition a harmonic measure (due to the integral definition of harmonicity), if it is conveniently normalized, and the corresponding Lyapunov exponent is negative. It gives us the desired a contradiction. The lemma is proved. $`\mathrm{}`$
To prove Theorem A in the case where $`=M`$, we consider the family of volume forms
$$\mu _n=\mathrm{exp}(q\psi _n)\mathrm{vol}_gv_t,$$
where the functions $`\psi _n`$ are given by Lemma 3.4. After normalizing them to probability measures and taking a subsequence, they converge to a probability measure $`\mu `$ on $`M`$.
###### Lemma 3.5.
The measure $`\mu `$ is transversely invariant.
Proof. It is more convenient to consider the family of transverse forms $`v_{t,n}=\mathrm{exp}(q\psi _n)v_t`$ as currents, i.e. as operators on the space of $`p`$-forms along the leaves. Define:
$$C_n(\omega )=_M\omega v_{t,n},$$
for every $`p`$-form $`\omega `$, where $`p=dim()`$. The family of currents $`\{C_n\}`$ is bounded, thus after taking a subsequence they converge to a current $`C`$. By construction, by choosing well the subsequence of currents converging to $`C`$, we have
$$\mu =\mathrm{vol}_gC.$$
We prove that $`C`$ is a closed current, or equivalently that $`C_n`$ is asymptotically closed as $`n`$ goes to infinity. Thus by Sullivan’s Theorem (\[Su1, Theorem I.12, p. 235\]) the measure $`\mu `$ is totally invariant.
Consider a $`(p1)`$-form $`\alpha `$. Using a partition of unity, we can write $`\alpha `$ as a finite sum of $`(p1)`$-forms whose support is contained in a foliation box. Thus we will suppose that the support of $`\alpha `$ is contained in a foliation box $`B\times T`$. Let $`\theta `$ be a volume form on $`T`$; write
$$v_{t,n}=\mathrm{exp}(\phi _n)\theta ,$$
then
(3.3)
$$\mu _n=v_{t,n}\mathrm{vol}_g=\mathrm{exp}(\phi _n)\theta \mathrm{vol}_g.$$
We have
$$_M𝑑\alpha v_{t,n}=_M𝑑\alpha \mathrm{exp}(\phi _n)\theta =$$
$$_M\alpha d(\mathrm{exp}(\phi _n)\theta )=_M(\alpha d\phi _n)v_{t,n}.$$
Thus, by Schwarz inequality,
(3.4)
$$\left|_M𝑑\alpha v_{t,n}\right|c|\alpha |_{\mathrm{}}_M|\phi _n|\mu _nc|\alpha |_{\mathrm{}}\left(_M|\phi _n|^2\mu _n\right)^{1/2},$$
where $`c`$ is a constant, and $``$ denotes the leafwise gradient. Observe that $`\phi _n`$ is well-defined up to addition of a leafwise constant function, so that $`\phi _n`$ is well-defined.
By Lemma 3.2:
$$_M|\phi _n|^2\mu _n=_M\mathrm{\Delta }\phi _n\mu _n,$$
and by Lemma 3.4, the Laplacians $`\mathrm{\Delta }\phi _n`$ verify uniformly
$$\underset{n\mathrm{}}{lim\; inf}\mathrm{\Delta }\phi _n0.$$
Thus the integrals
$$_M|\phi _n|^2\mu _n$$
tend to $`0`$ when $`n`$ goes to infinity. By formula 3.4, we conclude that $`C`$ is closed. Hence Lemma 3.5 is proved. $`\mathrm{}`$
The proof of this lemma concludes the proof of Theorem B when $`=M`$.
###### Remark 3.6.
In general, for a codimension $`q`$ foliation of class $`C^1`$, Oseledets’ Theorem states the existence of $`q`$ Lyapunov exponents $`\lambda _1(\mu )`$,…, $`\lambda _q(\mu )`$ associated to any harmonic ergodic measure $`\mu `$. In fact, the method we have used in this section proves that for any symmetric Laplace foliation of class $`C^1`$: either there exists a harmonic measure $`\mu `$ for which $`\lambda _1(\mu )+\mathrm{}+\lambda _q(\mu )`$ is negative, or there exists a transversely invariant measure. This result is analogue to the one of Baxendale \[Ba\].
However, let us mention some differences between the context of groups and the one of foliations. The theorem of Baxendale does not require the dynamics to be symmetric. It is interesting to note that Baxendale’s Theorem does not work for non symmetric Laplace foliations: there exists an example of a minimal foliation (drifted geodesic flow, see paragraph 3.4), for which every harmonic measure has positive Lyapunov exponent. The fact that the probability does not depend on the point in the Theorem of Baxendale should be interpreted in the foliation context by the concept of similarity. There is also an example of a similar Laplace lamination which does not verify Baxendale’s Theorem (in this setting the Lyapunov exponent should be defined by using the transverse $`2`$-adic structure).
#### 3.1.2. Proof of Theorem B in the exceptional minimal set case
First, let us notice that $``$ can not support a harmonic measure with positive Lyapunov exponent. A general idea implying this is the following. Consider first the case of a similar foliation. Then, a harmonic measure on $``$ induces a harmonic transverse measure (see Definition 4.3), which is harmonic. Due to the Ito Formula, the transversal measure of the image of a small transverse ball under the holonomy map $`h_{\gamma |_{[0,t]}}`$ associated to the Brownian path $`\gamma |_{[0,t]}`$ is a martingale.
Hence, its expectation at any Markovian moment equals the measure of the initial transverse ball. On the other hand, a small transverse ball around a typical point (in the sense of convergence of the Lyapunov exponent) is exponentially expanded by a typical Brownian path.
Thus, the martingale takes (at some big moments of time) large values with a large probability. Thus, the expectation is large. This contradicts the fact that the initial transverse ball can be chosen arbitrarily small. Rigorous proof of this statement is presented in Section 5, Lemma 5.4. Thus, all the Lyapunov exponents are equal to $`0`$.
The fact that the Laplace operator is symmetric was used when we stated that conditional measures are harmonic: if the Laplace operator is non-symmetric, these measures are harmonic in the sense of the adjoint operator $`\mathrm{\Delta }^{}`$, not $`\mathrm{\Delta }`$.
Now, let us continue the proof using the fact that all the Lyapunov exponents vanish. Choose some $`\epsilon >0`$. Then, the arguments of Lemma 3.4 imply that there exists a function $`\psi _\epsilon ^0`$, such that $`\epsilon <f\mathrm{\Delta }\psi _\epsilon ^0<\epsilon `$. By continuity, the same inequality holds in some neighborhood $`U^\epsilon `$ of $``$, that we suppose to be contained in the $`\epsilon `$-neighborhood of $``$.
###### Lemma 3.7.
If $`\epsilon >0`$ is small enough, there exists a function $`\psi _\epsilon `$, such that $`\mathrm{\Delta }\psi _\epsilon \epsilon `$ and such that at least ($`1\epsilon `$)-part of the measure $`\mu _\epsilon =e^{\psi _\epsilon ^0+\psi _\epsilon }\mathrm{vol}_gv_t`$ is concentrated in $`U^\epsilon `$.
Once such functions are constructed for any $`\epsilon `$, the proof will be finished in the same way as the proof in a minimal case. Namely, suppose that such functions are constructed. Let us find any weak limit $`\mu `$ of a subsequence of a family $`\frac{1}{\mu _\epsilon (M)}\mu _\epsilon `$ as $`\epsilon 0`$. Note that $`\mu `$ is supported on $``$: this comes from the fact that $`\mu _\epsilon (U^\epsilon )1\epsilon `$ and that $`_\epsilon U^\epsilon =`$. Also, note that $`\mu `$ is a weak limit for the same subsequence of the (non-normalized) measures $`\mu _\epsilon |_{U^\epsilon }`$. But these restricted measures can be written (locally) as $`e^\phi ^{}\mathrm{vol}_g\theta `$, where $`\theta `$ is a transverse measure, and $`\mathrm{\Delta }\phi ^{}=q\mathrm{\Delta }\phi q\mathrm{\Delta }\psi _\epsilon ^0+\mathrm{\Delta }\psi _\epsilon 2\epsilon `$. Passing to the limit and making estimates as in Lemma 3.5, we obtain that $`\mu `$ is totally invariant. Thus the proof of the Theorem will be complete after the proof of Lemma 3.7.
Proof of Lemma 3.7. Recall that $`\epsilon >0,\psi _\epsilon ^0`$ and $`U^\epsilon `$ are already chosen. We choose another neighborhood $`V`$ of $``$, $`VU^\epsilon `$. We are going to find a function $`\psi =\psi _\epsilon `$ verifying Lemma 3.7.
First, let us suppose that in $`U^\epsilon `$ there is no other minimal set, and thus that every leaf passing through a point in $`U^\epsilon `$ intersects $`U^\epsilon `$. We are going to look for the function $`\psi `$ as a solution of the Poisson equation
$$\mathrm{\Delta }\psi (x)=\epsilon ,xU^\epsilon V;\psi |_{(U^\epsilon V)}=0,$$
extended by $`0`$ to the complementary of $`U^\epsilon V`$. The solution of this problem always exists and can be found in the following way:
$$\psi (x)=\epsilon 𝔼T(\gamma ),$$
where $`T(\gamma )`$ denotes the first intersection moment of a Brownian path $`\gamma `$ with the boundary of $`U^\epsilon V)`$:
$$T(\gamma )=\mathrm{min}\{t:\gamma (t)(U^\epsilon V)\},$$
and $`𝔼`$ is the expectation of a function on the probability space $`(\mathrm{\Gamma }_x,W_x)`$. Note that $`\mathrm{\Delta }\psi `$ equals $`\epsilon `$ in $`U^\epsilon V`$, and $`0`$ in $`V`$ and in the complementary of $`U^\epsilon `$. Moreover, on $`(U^\epsilon V)`$, $`\mathrm{\Delta }\psi `$ is a positive distribution. Thus, one has $`\mathrm{\Delta }\psi \epsilon `$.
Now, let us show that for an appropriate choice of $`V`$ the major part of $`\mu _\epsilon `$ is concentrated in $`U^\epsilon `$, with the precise estimates of the Lemma. To do this, it suffices to check that as $`V`$ tends to $``$, the part of the measure $`\mu _\epsilon `$, concentrated in $`U^\epsilon `$, tends to 1. Note that the (non-normalized!) measure $`\mu _\epsilon `$ of $`MU^\epsilon `$ does not change, so that we have to prove that the measure of $`U^\epsilon `$ tends to infinity. By the monotone convergence theorem, it is equivalent to the fact that if we let $`\overline{\psi }=lim_V\psi `$ (maybe, $`\overline{\psi }`$ equals infinity at some points), the function $`e^{\overline{\psi }}`$ will be non-integrable in $`U^\epsilon `$. Note that the function $`\overline{\psi }`$ can be written as:
$$\overline{\psi }(x)=\epsilon 𝔼T_0(\gamma ),$$
where
$$T_0(\gamma )=\mathrm{min}\{t:\gamma (t)U^\epsilon \}$$
(if such an intersection does not occur, we define $`T_0(\gamma )=\mathrm{}`$). Thus, we have to estimate the mean $`𝔼T_0(\gamma )`$.
Note that in a neighborhood $`U^\epsilon `$ we have $`f\mathrm{\Delta }h_\epsilon <\epsilon `$. Thus, for a distance $`\stackrel{~}{d}`$ induced by a transversal metric $`e^{h_\epsilon }||`$, we have $`\mathrm{\Delta }\mathrm{log}\stackrel{~}{d}(,)<\epsilon `$ in $`U^\epsilon `$.
Now, let us consider a random process
$$\xi _0(t,\gamma )=\mathrm{log}\stackrel{~}{d}(\gamma (t),)\epsilon t,$$
and let us stop it at the moment $`T_0(\gamma )`$:
$$\xi (t,\gamma )=\xi _0(\mathrm{min}(t,T_0(\gamma )),\gamma ).$$
Then, the Ito formula implies that $`\xi (t,\gamma )`$ is a supermartingale:
$$\begin{array}{c}\frac{}{s}|_{s=t+0}𝔼(\xi (s,\gamma )|\gamma |_{[0,t]})=\hfill \\ \hfill =\{\begin{array}{cc}(\mathrm{\Delta }\mathrm{log}\stackrel{~}{d}(,))(\gamma (t))\epsilon ,\hfill & \gamma |_{[0,t]}U^\epsilon \hfill \\ 0,\hfill & T_0(\gamma )t.\hfill \end{array}0\end{array}$$
Note also, that the function $`\mathrm{log}\stackrel{~}{d}(,)`$ is Lipschitz on the leaves, thus, the conditional second moments
$$𝔼(r_n^2(\gamma )|\gamma |_{[0,t]})$$
of the increasements
$$r_n(\gamma )=\xi (n+1,\gamma )\xi (t,\gamma )$$
are bounded uniformly on $`n`$ and $`\gamma |_{[0,n]}`$. Thus, due to the theory of martingales, for every Markovian moment $`\tau `$ with finite expectation the expectation of $`\xi `$ at this moment does not exceed its initial value.
Let us now use this process to estimate from below the expectation $`𝔼T_0(\gamma )`$. Either this expectation is infinite (in which any lower bound is satisfied automatically). Or it is finite, and in this case the expectation of a value of a supermartinagle $`\xi `$ in a Markovian moment $`T_0(\gamma )`$ does not exceed its initial value, that is
$$𝔼\left[\mathrm{log}\stackrel{~}{d}(\gamma (T_0(\gamma )),)\epsilon T_0(\gamma )\right]\mathrm{log}\stackrel{~}{d}(x,).$$
The expectation in the left side can be rewritten as
$$\epsilon 𝔼T_0(\gamma )+𝔼\mathrm{log}\stackrel{~}{d}(\gamma (T_0(\gamma )),)=\epsilon 𝔼T_0(\gamma )+O(1),$$
for at the moment of exiting $`U^\epsilon `$ the distance to $``$ is separated from $`0`$. So, we have
$$\overline{\psi }=𝔼T_0(\gamma )\frac{1}{\epsilon }\mathrm{log}\stackrel{~}{d}(x,)+C_0.$$
This implies that
$$e^{\overline{\psi }(x)}\frac{C}{(\stackrel{~}{d}(x,))^{1/\epsilon }}.$$
Thus, if $`\epsilon `$ is less than $`1/(codim)`$, the function $`e^{\overline{\psi }}`$ is non-integrable and the effect of concentration takes place. This completes the proof under the hypothesis that $``$ is the unique minimal subset of $`U^\epsilon `$.
To conclude the proof, we remark that if $`U^\epsilon `$ contains another minimal set, we can replace $``$ by the closure of the union of all the leaves, entirely contained in $`U^\epsilon `$, and repeat the previous arguments. $`\mathrm{}`$
The proof of Theorem B is completed in all the cases.
### 3.2. Proof of the Main Theorem
We suppose that the foliation $``$ is transversely conformal and does not have a transversely invariant measure. By Theorem B, on any minimal set is supported a unique harmonic measure with negative Lyapunov exponent. Because of the attraction property, any minimal set has a neighborhood which does not contain any other minimal set. Thus, there is a finite number of minimal sets $`_1,\mathrm{},_k`$. Denote by $`\mu _1,\mathrm{},\mu _k`$ their unique harmonic measure, and $`\lambda _1,\mathrm{},\lambda _k`$ the corresponding Lyapunov exponents. Note that every point $`xM`$ belongs to the basin of attraction of at least one of these sets; the reason is that the set
$$M(\underset{j=1}{\overset{k}{}}Attr(_j))$$
is closed, consists only of entire leaves and does not contain any minimal subset.
Let $`\alpha >0`$ be a real number such that $`\alpha <|\lambda _j|`$ for every $`j`$. For every point $`xM`$ we consider the probability
$$p_j(x)=W_x(\{\gamma \mathrm{\Gamma }_x\gamma (t)\underset{t\mathrm{}}{\overset{}{}}_j\}),$$
that a Brownian path starting at $`x`$ tends to $`_j`$. Note that almost every Brownian path tending to $`_j`$ (if such path exists) is distributed with respect to $`\mu _j`$, and contracts a transverse ball at $`x`$ exponentially with exponent $`\alpha `$ (see Remark 2.8). We claim that the sum of these probabilities is equal to $`1`$; in other words, $`W_x`$-almost every trajectory tends to one of the minimal sets, with the distribution and transverse contraction properties. We show this in the following way: for arbitrary small neighborhoods $`U_1,\mathrm{},U_k`$ of $`_1,\mathrm{},_k`$ respectively, the complementary $`R=M_jU_j`$ is a closed set without any minimal subset, thus containing no entire leaf. Hence, for any point $`x`$ of $`R`$ there exists a leafwise path leading to one of the neighborhoods $`U_j`$; moreover, by compactness of $`M`$, the length of such a path is bounded uniformly on $`R`$. Thus, for a point $`xR`$, the probability that it lies in one of the $`U_j`$ at time $`1`$ is bounded from below by a positive uniform constant. Hence, for any point $`xM`$, almost every trajectory starting in $`x`$ meets one of the neighborhoods $`U_j`$.
To complete the proof, let us show that $`_jp_j(x)>1\epsilon `$ for any $`\epsilon >0`$. To do this, let us choose $`U_j`$ so close to $`_j`$ that for every point in $`U_j`$ the probability of attracting to $`_j`$ with the distribution and the transverse contraction properties is at least $`1\epsilon `$; it is possible due to Lemma 2.4. Now, let us use the Markovian property: for any $`xM`$, almost every trajectory $`\gamma \mathrm{\Gamma }_x`$ meets one of the $`U_j`$, and for a starting point in $`U_j`$ the probability of attracting to the corresponding $`_j`$ is at least $`1\epsilon `$. Thus, the probability of attracting to one of the $`_j`$ is at least $`1\epsilon `$. As $`\epsilon >0`$ was chosen arbitrary, we have proven that almost every trajectory tends to one of the $`U_j`$.
Recall that the functions $`p_j`$ are leafwise harmonic and lower semicontinuous. Because their sum is equal to the constant function $`1`$ the functions $`p_j`$ are continuous.
Thus, we have proved the Contraction, Distribution and Attraction parts of the main theorem. We end the proof of Theorem 1.1 by proving the statement about the asymptotic behaviour of the diffusion. First, we shall prove a weaker form. Namely, we prove that the time-averages of diffusions tends to the same limit:
$$\frac{1}{T}_0^TD^tf𝑑t\stackrel{xM}{}\psi (x),$$
where $`\psi (x)=_jp_jf𝑑\mu _j`$. In the case where the foliation is minimal, it is implied by unique ergodicity (see analogous arguments in \[Fu2\]). Namely, the value of the time-average of the diffusions at a point $`xM`$ can be rewritten as an integral:
$$\begin{array}{c}\frac{1}{T}_0^T(D^tf)(x)𝑑t=\frac{1}{T}_0^T_MD^tf𝑑\delta _x𝑑t=\hfill \\ \hfill =\frac{1}{T}_0^T_Mfd(D_{}^t\delta _x)𝑑t=_Mf𝑑m_{x,T},\end{array}$$
where
$$m_{x,T}=\frac{1}{T}_0^T(D_{}^t\delta _x)𝑑t$$
is the time-average of the diffusions of the measure $`\delta _x`$. Note that due to a classical argument in ergodic theory, a weak limit of a sequence $`m_{x_n,t_n}`$ with $`t_n\mathrm{}`$ is harmonic. As there exists a unique harmonic measure $`\mu `$, the time averages $`m_{x,t}`$ converge to $`\mu `$ uniformly in $`x`$ as $`t`$ tends to infinity. Thus, the integrals of $`f`$ with respect to these measures also converge uniformly to $`_Mf𝑑\mu `$, which implies the desired statement.
In the case of an exceptional minimal set we notice that the time-averages of the diffusions can be rewritten as
(3.5)
$$\frac{1}{T}_0^T(D^tf)(x)𝑑t=_{\mathrm{\Gamma }_x}\left(\frac{1}{T}_0^Tf(\gamma (t))𝑑t\right)𝑑W_x(\gamma ).$$
We know that $`W_x`$-almost all trajectories tend to one of the $`_j`$’s and are distributed with respect to the corresponding harmonic measure $`\mu _j`$. The probability that a point $`x`$ tends to $`_j`$ is equal to $`p_j(x)`$. Hence, the right hand side of (3.5) is equal to
$$\underset{j=1}{\overset{k}{}}p_j(x)__jf𝑑\mu _j.$$
Moreover, a uniform argument on $`(1\epsilon )`$-measure of trajectories (similarity of Brownian motions) implies that this convergence is uniform in $`x`$.
Thus, in every case we have shown that the time-average of the diffusions converge to the right-hand side, which we denote $`\psi `$.
Now, let us finish the proof using the arguments analogous to these of Kaimanovich \[Kai1\]. Namely, notice that due to the diffusion of the Brownian motion, there exists $`\epsilon ,\epsilon ^{}>0`$ such that for any point $`x`$, and any point $`yU_\epsilon (x)`$, the densities $`p(x,y,1)`$ and $`p(x,y,2)`$ are bounded from below by $`\epsilon ^{}`$. Thus one has for every $`x`$ and every bounded function $`f`$:
$$|D^1f(x)D^2f(x)|2(1\epsilon ^{}\mathrm{vol}(U_\epsilon )(x))|f|_{\mathrm{}},$$
where $`||_{\mathrm{}}`$ is the uniform norm. Thus, because the leaves are of bounded geometry, we have
$$D^1D^2_{\mathrm{}}<2,$$
where $`||||_{\mathrm{}}`$ is the norm of operators acting on $`L^{\mathrm{}}`$. The “zero-two law” \[Li\] implies that $`D^nD^{n+1}0`$ as $`n\mathrm{}`$. In particular, the time-averages of the diffusions converge if and only if the diffusions converge themselves to the same limit \[Kai2\]. Hence, the diffusions converge to the limit we have described.
Proof of Corollary 1.3. Let $``$ be a transversely conformal foliation of class $`C^1`$ of a compact manifold, and $``$ a minimal set of $``$. By Theorem B, either $``$ supports a transversely invariant measure, or a harmonic measure of negative Lyapunov exponent. In this case, Candel has proved that there exists a loop contained in a leaf of $``$ with hyperbolic holonomy (see \[Can, Theorem 8.18\]). $`\mathrm{}`$
### 3.3. Examples: codimension one foliations of class $`C^2`$
In the case of codimension one foliations of class $`C^2`$ without compact leaf, the following result completes the Main Theorem:
###### Proposition 3.8.
Let $``$ be a codimension one foliation of class $`C^2`$ of a compact manifold, without compact leaves. Then if $``$ has a totally invariant measure $`\mu `$, this measure is the unique harmonic measure. Then, for every point $`x`$, almost every Brownian path starting at $`x`$ is distributed with respect to $`\mu `$, and the diffusions of a continuous function $`f:M`$ tend uniformly to the constant function $`f𝑑\mu `$.
Proof. By Sacksteder Theorem, the foliation $``$ is minimal. We first prove that the measure $`\mu `$ is the unique totally invariant measure. By minimality of $``$ and Haefliger’s argument \[Hae\], there exists a transverse circle $`C`$ cutting every leaf. The transversely invariant measure (corresponding to $`\mu `$) induces a measure $`\theta `$ on $`C`$, invariant by all the holonomy maps. This measure gives us a map $`h:C/l=C^{}`$, where $`l=\theta (C)`$. This map semi-conjugates the pseudo-group induced by $``$ on $`C`$ to a finitely generated group of rotations of $`C^{}`$, which we note $`G`$. Because $``$ is minimal, and the holonomy pseudo-group is finitely generated, at least one of the rotations of $`G`$ is irrational. Then, the Lebesgue measure is the unique probability measure invariant by $`G`$, and thus $`\mu `$ is the unique totally invariant measure on $``$ up to multiplication by a constant.
To conclude the case when there exists a totally invariant measure $`\mu `$, it suffices to show that every harmonic measure is in fact a totally invariant measure. Observe that the group $`G`$ is a group of rotations, so that the orbits of its action on $`C^{}`$ have a polynomial growth. Hence, the same is true for the action of the holonomy group on $`C`$. Thus, every leaf grows polynomially. Kaimanovich (\[Kai1, Corollary of Theorem 4\]) proved that if for a harmonic measure, almost every leaf (with respect to this measure) has subexponential growth, then this measure is totally invariant. In our case, all the leaves have polynomial growth, hence every harmonic measure is in fact totally invariant.
The distribution and diffusion property is implied by the fact that the harmonic measure is unique. $`\mathrm{}`$
We end the paragraph by constructing a foliation by surfaces of a $`3`$-dimensional compact manifold, with two exceptional minimal sets. The example is constructed in the following way. Let $``$ be an oriented codimension one foliation by oriented surfaces with an exceptional minimal set $``$ and suppose that there exists a transverse loop $`c`$ which does not cut $``$. Then a neighborhood of $`c`$ in $`M`$ is diffeomorphic to a solid torus $`D^2\times 𝕊^1`$, the foliation $``$ being the horizontal fibration by two dimensional balls $`D^2`$. Now consider two copies $`N_1`$ and $`N_2`$ of the exterior of $`D^2\times 𝕊^1`$ in $`M`$. These two manifolds are foliated, and have a boundary component $`B\times 𝕊^1`$ transverse to the foliation $``$. The foliation $``$ induces the horizontal foliation by circles on it. Observe also that $`N_1`$ and $`N_2`$ have an exceptional minimal set in their interior. Thus, by gluing $`N_1`$ and $`N_2`$ along their boundary by a diffeomorphism which preserves the foliation and reverses the orientation, we construct a foliation by surfaces of a closed manifold with two exceptional minimal sets.
Now we are given an example of such a situation. We consider a surface $`\mathrm{\Sigma }`$ of bounded topology and constant negative curvature, with a cusp of infinite volume. The cusp determines an interval $`I`$ in the boundary of the universal cover of $`\mathrm{\Sigma }`$ which has two remarcable properties. The first is that it is invariant by the action of the geodesic $`\gamma `$ on $`\stackrel{~}{\mathrm{\Sigma }}`$. The second is that it is a component of the exterior of the limit set of $`\pi _1(\mathrm{\Sigma })`$. Now consider a compact surface $`S`$ of sufficiently large genus so that there exists a surjective morphism $`\rho :\pi _1(S)\pi _1(\mathrm{\Sigma })`$. We get an action of the fundamental group of $`S`$ on the boundary of $`\stackrel{~}{\mathrm{\Sigma }}`$, which leaves the limit set of $`\pi _1(\mathrm{\Sigma })`$ invariant, and for which there exists an element leaving $`I`$ invariant, and which acts as a translation on it. Let $`(M,)`$ be the supension of $`\rho `$: this is the foliation induced by a flat circle bundle over $`S`$ whose holonomy is smoothly conjugated to the representation $`\rho `$ (see 4.1.1). Then the saturated subset of the limit set of $`\pi _1(\mathrm{\Sigma })`$ is an exceptional minimal subset $``$ of $``$. Now, by construction, there is a leaf $`L`$ which intersects twice the component $`I`$ of the exterior of the limit set of $`\pi _1(\mathrm{\Sigma })`$ in $`\stackrel{~}{\mathrm{\Sigma }}`$. By the standard Haefliger’s argument, we construct a transverse circle which does not cut $``$. Thus applying the preceeding arguments we construct a foliation with two minimal sets.
### 3.4. A counter-example in the non symmetric case
In \[Can\], Candel extends Garnett’s theory to the case of non symmetric Laplace operators on a foliation. In the case of non symmetric Laplace operators, the dichotomy “The Lyapunov exponent is positive or there exists a transversely invariant measure” does not hold anymore. In this paragraph we describe a nice counter-example in the non symmetric case (see also \[Ham2\]).
Consider a compact Riemannian manifold $`(M,g)`$ of dimension $`3`$ on which there is an orthonormal frame $`(H^s,V,H^u)`$, for which the vector fields $`H^s,V,H^u`$ verify the relations:
$$[V,H^s]=H^s,[V,H^u]=H^u,[H^u,H^s]=V.$$
Such manifolds are quotient of the universal cover of $`SL(2,𝐑)`$ by a cocompact lattice $`\mathrm{\Gamma }`$. If $`\mathrm{\Sigma }`$ is the quotient of the upper-half plane $`𝐇`$ by $`\mathrm{\Gamma }`$, then $`M`$ is naturally identified with the unitary tangent bundle of $`\mathrm{\Sigma }`$, and under this identification $`V`$ is the geodesic flow of the hyperbolic surface, $`H^s`$ and $`H^u`$ the horocycle flows. The vector fields $`V`$ and $`H^s`$ generate a foliation $`^s`$ which is the stable foliation of the flow $`V`$.
Let $`g^s`$ be the restriction of the metric $`g`$ on $`^s`$. For any $`\kappa `$, consider the Laplace operator $`\mathrm{\Delta }_\kappa `$ defined by
$$\mathrm{\Delta }_\kappa =\mathrm{\Delta }_{g^s}+\kappa V,$$
where $`\mathrm{\Delta }_{g^s}`$ is the Laplacian of $`g^s`$ along the leaves of $`^s`$. Garnett proved that for the symmetric case $`\kappa =0`$, the Liouville measure $`\mathrm{vol}_g`$ is the unique harmonic measure (see \[Ga, Proposition 5, p. 305\]). Note that the Liouville measure on $`M`$ is also invariant by $`V`$, so that it is a harmonic measure for all the Laplace operators $`\mathrm{\Delta }_\kappa `$.
###### Theorem 3.9.
For any $`\kappa `$, the Lyapunov exponent of any harmonic measure $`\mu `$ of $`(^s,\mathrm{\Delta }_\kappa )`$ is $`\lambda (\mu )=\kappa 1`$. When $`\kappa <1`$ the Liouville measure is the unique harmonic measure. When $`\kappa >1`$, there exists a harmonic measure supported on every cylinder leaf (thus the foliation is not uniquely ergodic).
Proof. First, we compute the Lyapunov exponent of a harmonic measure $`\mu `$ of $`(^s,\mathrm{\Delta }_\kappa )`$. To this end we use the formula of Lemma 3.1:
$$\lambda (\mu )=_M\mathrm{\Delta }_\kappa \mathrm{log}\phi d\mu .$$
Consider the metric $`||`$ on the normal bundle of $``$ which is induced by $`g`$. We are going to compute the function $`\phi `$, which is defined up to multiplication by a constant. This function verifies the relations
$$V\phi =\phi ,H^s\phi =0.$$
In the leaves we have local coordinates $`z=x+iy`$ with values in the upper half-plane $`𝐇`$, such that
$$V=y\frac{}{y},H^s=y\frac{}{x}.$$
(These coordinates are well defined up to an affine transformation of the upper half-plane). In these coordinates, $`\phi =y`$ up to a multiplicative constant. The metric $`g^s`$ and the Laplacian $`\mathrm{\Delta }_{g^s}`$ are expressed by
$$g^s=\frac{dx^2+dy^2}{y^2},\mathrm{\Delta }_{g^s}=y^2(\frac{^2}{x^2}+\frac{^2}{y^2}).$$
Thus, we have $`\mathrm{\Delta }_\kappa \phi =\kappa 1`$ identically, and the formula $`\lambda (\mu )=\kappa 1`$ follows. In particular, when $`\kappa <1`$ the only harmonic measure is the Liouville measure, because of Theorem A.
Now, let us suppose that $`\kappa >1`$. We shall prove that every cylinder leaf supports a harmonic measure. Observe that these leaves are those containing a periodic orbit of the vector field $`V`$. Let $`L`$ be such leaf, $`\gamma _0`$ be the closed orbit of $`V`$ in $`L`$. Then, its universal cover is the hyperbolic plane, for which we choose the upper half-plane model $`𝐇`$. Observe that we have the canonical coordinates up to an affine transformation, constructed before. Without loss of generality, we may suppose that the geodesic in $`𝐇`$, corresponding to $`\gamma _0`$, is the vertical geodesic $`x=0`$ going upwards. Denote by $`A`$ the length of $`\gamma _0`$; then the transformation of $`𝐇`$ corresponding to $`\gamma _0`$ as an element of $`\pi _1(L)`$, is $`ze^Az`$. Thus, the leaf $`L`$ is obtained from $`𝐇`$ by identifying $`z`$ and $`e^Az`$.
Now, let us consider a typical Brownian trajectory $`\gamma `$ in $`L`$ and its lift to the universal cover $`\stackrel{~}{\gamma }(t)=(x(t),y(t))`$. Note that $`\stackrel{~}{\gamma }`$ satisfies the following stochastic differential equation:
$$\{\begin{array}{c}\dot{x}=\sqrt{2}ydW_t^1,\hfill \\ \dot{y}=cy+\sqrt{2}ydW_t^2,\hfill \end{array}$$
where $`W_t^1`$ and $`W_t^2`$ are two independent Wiener processes. Here the coefficient $`\sqrt{2}`$ comes from our definition of Brownian motion: as we have defined it using the heat kernel, its intensivity equals $`2`$ (instead of its common value $`1`$). Let us make a change of variables: let $`u=\mathrm{log}y`$, $`v=x/y`$. Then
$$\begin{array}{c}\dot{u}=(\mathrm{log}y)^{}=(\mathrm{log})^{}(y)\kappa y+\frac{1}{2}(\mathrm{log})^{\prime \prime }(y)2y^2+\hfill \\ \hfill +(\mathrm{log})^{}(y)\sqrt{2}ydW_t^1=(1+\kappa )+\sqrt{2}dW_t^1,\end{array}$$
$$\begin{array}{c}\dot{v}=(\frac{x}{y})_y^{}\kappa y+\frac{1}{2}(\frac{x}{y})_{yy}^{\prime \prime }2y^2+(\frac{x}{y})_x^{}\sqrt{2}ydW_t^1+(\frac{x}{y})_y^{}\sqrt{2}ydW_t^2=\hfill \\ \hfill =\kappa y\frac{x}{y^2}+\frac{2x}{y^3}y^2+\sqrt{2}(dW_t^1+\frac{x}{y}dW_t^2)=(2\kappa )v+\sqrt{2}(dW_t^1+vdW_t^2).\end{array}$$
This implies that $`v`$ satisfies a stochastic differential equation
$$\dot{v}=(2\kappa )v+\sqrt{2(1+v^2)}dW_t.$$
Let us now make another change of variable: we denote $`\xi =f(v)=\mathrm{log}(v+\sqrt{1+v^2})`$. Then $`\xi `$ satisfies the following stochastic differential equation:
$$\begin{array}{c}\dot{\xi }=f^{}(\xi )(2\kappa )v+\frac{1}{2}f^{\prime \prime }(\xi )(\sqrt{2(1+v^2)})^2+f^{}(\xi )\sqrt{2(1+v^2)}dW_t=\hfill \\ \hfill =\frac{1}{\sqrt{1+v^2}}(2\kappa )v\frac{1}{4}\frac{2v}{\sqrt{1+v^2}^3}(2(1+v^2))+\sqrt{2}dW_t=\\ \hfill =\frac{v}{\sqrt{1+v^2}}(1\kappa )+\sqrt{2}dW_t.\end{array}$$
For $`\kappa >1`$ we notice that the Brownian component of this stochastic differential equation is constant, and the drift is towards $`0`$ with the velocity separated from zero for large $`v`$. Thus, there exists a probability stationary measure for this process on the real line. By lifting this measure to the initial cylinder (by a product with the Lebesgue measure in $`\mathrm{log}y`$), we obtain a stationary measure on $`L`$. We have constructed a harmonic measure on $`L`$. $`\mathrm{}`$
## 4. Similar foliations
A Laplace foliation $`(,\mathrm{\Delta })`$ is called similar if there exists a transverse continuous foliation $`𝒢`$ of dimension $`codim()`$ such that the operator $`\mathrm{\Delta }`$ is invariant by $`𝒢`$; thus $`𝒢`$ preserves the metric $`g`$ and the drift vector field $`V`$.
The main goal of this part is to prove the unique ergodicity property for a codimension $`1`$ similar Laplace foliation whose drift vector field preserves the volume, and whose transverse structure is just supposed continuous. We begin by giving examples of such foliations.
### 4.1. Some examples
#### 4.1.1. Suspension
Let $`(N,\overline{\mathrm{\Delta }})`$ be a compact manifold equipped with a Laplace operator $`\overline{\mathrm{\Delta }}`$, and $`\rho :\pi _1(N)Homeo(F)`$ be a representation of its fundamental group into the group of homeomorphisms of a compact manifold $`F`$. Let $`\stackrel{~}{N}`$ be the universal cover of $`N`$. The diagonal action of the discrete group $`\pi _1(N)`$ on the product $`\stackrel{~}{N}\times F`$ is discontinuous and free. Moreover, it preserves the horizontal foliation and the vertical fibration. Thus, the quotient $`N_\rho F`$ is equipped with a foliation $``$ (quotient of the horizontal foliation) and with a transverse fibration $`FM\stackrel{\pi }{}N`$ (quotient of the vertical fibration). Let $`\mathrm{\Delta }`$ be the Laplace operator on the leaves of $``$ such that $`(\pi _{})_{}\mathrm{\Delta }=\overline{\mathrm{\Delta }}`$. By construction the foliation $`(,\mathrm{\Delta })`$ is similar. Such foliations are called suspensions.
#### 4.1.2. Linear Anosov diffeomorphism
Let $`A:𝐓^n𝐓^n`$ be a linear Anosov diffeomorphism of the torus $`𝐓^n=𝐑^n/𝐙^n`$. Consider the quotient $`M`$ of $`(0,\mathrm{})\times 𝐓^n`$ by the diffeomorphism $`\stackrel{~}{A}(t,x)=(2t,Ax)`$: this is the fiber bundle over the circle whose fiber is $`𝐓^n`$ and monodromy is given by $`A`$. Define the foliations $``$ and $`𝒢`$ to be respectively the quotient of $`(0,\mathrm{})\times ^u`$ and of $`^s`$. We define a Laplace operator $`\mathrm{\Delta }`$ on the leaves of $``$ of the form
$$\mathrm{\Delta }=\mathrm{\Delta }_t+t^2\frac{^2}{t^2},$$
where $`\{\mathrm{\Delta }_t\}_{t>0}`$ is a family of linear Laplace operators on $`^u`$ depending smoothly on $`t`$, and verifying the relation
$$\mathrm{\Delta }_{2t}=(A|_^u)_{}\mathrm{\Delta }_t,$$
for every $`t>0`$. The Laplace foliations $`(,g)`$ are similar (the invariance of $`\mathrm{\Delta }`$ by $`𝒢`$ comes from the fact that the operators $`\mathrm{\Delta }_t`$ are linear).
###### Remark 4.1.
The similar foliations by surfaces of a compact $`3`$-manifold are well known \[Ca, Ep\]: in this case, the only examples with interesting dynamics are the suspensions and the foliations induced by a linear Anosov diffeomorphism of a $`2`$-torus. However, it may exist other examples in higher dimension.
### 4.2. Non divergence of the leaves
In this paragraph we prove that the leaves of a similar codimension $`1`$ foliation whose drift vector field preserves the volume $`\mathrm{vol}_g`$ are not diverging in a set of directions of large measure. This has been observed by Thurston (see \[C-D, Fe\] for a topological proof).
###### Definition 4.2.
Let $``$ be a similar foliation. A transverse harmonic measure on $``$ is a family $`\{\nu _L\}`$ of measures $`\nu _L`$ on every $`𝒢`$-leaf $`L`$, such that in a chart $`B\times T`$ in which $``$ and $`𝒢`$ are respectively the horizontal and vertical foliations, the function
$$pB\nu (\{p\}\times T)𝐑_+$$
is harmonic.
###### Lemma 4.3.
On a codimension one similar Laplace foliation of a compact manifold whose drift vector field preserves the leafwise volume $`\mathrm{vol}_g`$, and which is minimal, there exists a transverse harmonic measure.
Proof. The result follows from the existence of a harmonic measure for the adjoint operator $`\mathrm{\Delta }^{}`$ of $`\mathrm{\Delta }`$. Recall that $`\mathrm{\Delta }^{}`$ is defined on every $``$-leaf $`L`$ in such a way that for any smooth functions $`u,v:L𝐑`$ with compact support one has
$$_Lu\mathrm{\Delta }v𝑑\mathrm{vol}_g=_L(\mathrm{\Delta }^{}u)v𝑑\mathrm{vol}_g.$$
An integration by parts shows that one has the following formula:
$$\mathrm{\Delta }^{}=\mathrm{\Delta }_gV+\mathrm{div}_{\mathrm{vol}_g}V,$$
where $`V`$ is the drift vector fields of $`\mathrm{\Delta }`$ (i.e. by definition $`\mathrm{\Delta }=\mathrm{\Delta }_g+V`$). If $`V`$ preserves the volume $`\mathrm{vol}_g`$, which means that the divergence of $`V`$ vanishes identically, then the operator $`\mathrm{\Delta }^{}`$ is also a Laplace operator. In \[Can\], it is proved that for such operators there exists a harmonic measure.
Let $`\mu `$ be a harmonic measure on $`(,\mathrm{\Delta }^{})`$. Consider a foliation box $`B\times T`$ in which $``$ and $`𝒢`$ are respectively the horizontal and vertical foliation. Let $`T^{}`$ be an open subset of $`T`$. The image of the measure $`\mu `$ by the projection $`B\times T^{}B`$ is a $`\mathrm{\Delta }^{}`$-harmonic measure on $`B`$, because the foliation $``$ is similar. Thus, there is a $`\mathrm{\Delta }`$-harmonic function $`L_T^{}:B[0,\mathrm{})`$ such that for any continuous function $`fC_c^0(B)`$, one has $`_{B\times T^{}}f(b)𝑑\mu (b,t)=L_T^{}(b)f(b)\mathrm{vol}_g(b)`$. Because $`\mu `$ is a measure, if $`T_n`$ are disjoint open subsets of $`T`$, one has the relations $`_nL_{T_n}=L_{_nT_n}`$; thus there exists a transverse measure on the leaves of $`𝒢`$ such that $`\nu (b\times T^{}):=L_T^{}(b)`$, for every Borel subset $`T^{}`$ of $`T`$. This transverse measure is $`\mathrm{\Delta }`$-harmonic, by construction and the lemma is proved. $`\mathrm{}`$
###### Lemma 4.4.
Let $``$ be a codimension one similar Laplace foliation of a compact manifold. Let $`[x,y]`$ be an interval in a $`𝒢`$-orbit. There exists a uniquely defined map $`I_{x,z}:\stackrel{~}{L_x}\times [x,y]M`$ which maps the horizontal Laplace foliation (given by the Laplace operator on $`\stackrel{~}{L_x}`$) on $`(,\mathrm{\Delta })`$, the vertical foliation on $`𝒢`$, and the interval $`\{x\}\times [x,y]`$ identically on $`[x,y]`$.
Proof. Let $`\gamma :[0,1]\stackrel{~}{L_x}`$ be a smooth path starting at $`x`$. The restriction of $`I_{x,y}`$ to $`\gamma ([0,1])\times [x,y]`$ is uniquely defined, if it exists. Let $`0t1`$ be the supremum of those $`t`$ such that $`I_{x,y}`$ is defined on $`\gamma ([0,t])\times [x,y]`$. It is clear that $`t>0`$, because locally we have foliation charts. Recall that the foliation $`𝒢`$ preserves the metric $`g`$; thus, for any $`z[x,y]`$ the length of the curve $`I_{x,y}(\gamma ([0,s]\times z)`$ equals the length of the curve $`\gamma ([0,s])`$. This implies that it is possible to extend the map $`I_{x,y}`$ on the domain $`\gamma ([0,t])\times [x,y]`$, and because locally we have foliation charts, to a domain $`\gamma ([0,t_+))\times [x,y]`$, where $`t_+>t`$. Thus $`t=1`$ and the lemma is proved. $`\mathrm{}`$
Recall that for any point $`xM`$, $`\mathrm{\Gamma }_x`$ is the set of continuous paths contained in the leaf through the point $`x`$, and the Laplace operator $`\mathrm{\Delta }`$ induces a probability measure $`W_x`$ on $`\mathrm{\Gamma }_x`$. Here is the main result of the paragraph, where $`J_{x,y}(p)`$ is the point $`I_{x,y}(p,1)`$.
###### Proposition 4.5.
Let $``$ be a Laplace similar foliation of codimension $`1`$ of a compact manifold, which is minimal, and whose drift vector field preserves the volume $`\mathrm{vol}_g`$. For every $`\epsilon >0`$, there exists a constant $`\delta >0`$ such that if $`x`$ and $`y`$ are two points on the same $`𝒢`$-orbit with $`d(x,y)\delta `$, then there exists a subset $`E_{x,y}\mathrm{\Gamma }_x`$ of $`W_x`$-measure $`1/2`$ such that for any $`\gamma E_{x,y}`$, one has
$$\underset{t\mathrm{}}{lim\; sup}d(\gamma (t),J_{x,y}(\gamma (t))\epsilon .$$
Proof. Consider a transverse harmonic measure $`\nu `$ constructed in Lemma 4.3. Because $``$ is minimal, the measure $`\nu `$ restricted to every leaf of $`𝒢`$ has full support, and is diffuse. Thus, there exists $`\delta ^{}>0`$ such that if $`[x,y]`$ is an interval in a $`𝒢`$-leaf of $`\nu `$-measure bounded by $`\delta ^{}`$, then the distance between $`x`$ and $`y`$ is bounded by $`\epsilon `$. Moreover, there exists $`\delta >0`$ such that if the distance between $`x`$ and $`y`$ is bounded by $`\delta `$, then the $`\nu `$-measure of the interval $`[x,y]`$ is bounded by $`\delta ^{}/2`$.
The function $`p\stackrel{~}{L_x}f(\gamma )=\nu (I_{x,y}(\{p\}\times [x,y]))(0,\mathrm{})`$ is harmonic. Thus by the martingale theorem, for $`w_x`$-almost every $`\gamma `$, the limit
$$\underset{t\mathrm{}}{lim}f(\gamma (t))$$
exists, and its integral over $`\mathrm{\Gamma }_x`$ is $`f(x)=\nu ([x,y])\delta ^{}/2`$. In particular, there is a measurable subset $`E_{x,y}\mathrm{\Gamma }_x`$ of $`W_x`$-measure $`1/2`$ such that for every $`\gamma E_{x,y}`$
$$\underset{t\mathrm{}}{lim}f(\gamma (t))\delta ^{}.$$
For every $`\gamma `$ of $`E_{x,y}`$ we have
$$\underset{t\mathrm{}}{lim\; sup}d(\gamma (t),I_{x,y}(\gamma )(t,1))\epsilon .$$
The proposition is proved. $`\mathrm{}`$
### 4.3. Application to unique ergodicity
We prove that there is only one harmonic measure on a codimension $`1`$ similar foliation whose drift vector field preserves the volume $`\mathrm{vol}_g`$.
###### Theorem 4.6.
Let $`(,\mathrm{\Delta })`$ be a similar Laplace foliation of codimension $`1`$ of a compact manifold $`M`$, whose drift vector field preserves the volume $`\mathrm{vol}_g`$. Then, on a minimal subset of $``$ is supported a unique harmonic measure.
Proof. Let $``$ be a minimal subset of $``$. There are three possibilities:
* $``$ is a compact leaf.
* $``$ is transversely a Cantor set.
* $`=M`$.
In the first case, the only harmonic measure is the unique $`\mathrm{\Delta }`$-harmonic volume on the compact leaf. The second case can be reduced to the third one by collapsing the components of the leaf of $`𝒢`$ outside $``$. Thus, we suppose that $`=M`$, i.e. $``$ is minimal.
###### Lemma 4.7.
Let $``$ be a similar minimal foliation of codimension $`1`$ of a compact manifold $`M`$, and $`\mu `$ be an ergodic harmonic measure. Then for every $`\alpha >0`$, and every continuous function $`f:M𝐑`$, the following property holds. For every point $`yM`$, there exists a measurable subset $`E_y\mathrm{\Gamma }_y`$ of $`W_y`$-measure $`1/2`$ such that for every $`\gamma E_y`$:
$$f𝑑\mu \alpha \underset{n\mathrm{}}{lim\; inf}B_n(f,\gamma )\underset{n\mathrm{}}{lim\; sup}B_n(f,\gamma )f𝑑\mu +\alpha ,$$
where $`B_n(f,\gamma ):=\frac{1}{n}_{1kn}f(\gamma (k))`$ are the Birkhoff sums of $`f`$ along the path $`\gamma `$.
Proof. By the Birkhoff theorem for harmonic measures proved in \[Ga\], there is a measurable subset $`XM`$ of full $`\mu `$-measure, saturated by $``$, so that for every $`xX`$ and $`w_x`$-almost every continuous path $`\gamma \mathrm{\Gamma }_x`$, the Birkhoff sums $`B_n(f,\gamma )`$ converge to $`f𝑑\mu `$.
Let $`\epsilon >0`$ such that if $`d(x,y)\epsilon `$ then $`|f(x)f(y)|\alpha `$. Let $`y`$ be any point of $`M`$. There exists a point $`xX`$ which is in the $`𝒢`$-orbit of $`y`$ and such that $`d(x,y)\delta `$, the $`\delta `$ being given by Lemma 4.5. Let $`F_x`$ denote the set of element $`\gamma \mathrm{\Gamma }_x`$ for which the Birkhoff sums converge to $`f𝑑\mu `$. Because $`x`$ belongs to $`X`$ this set is of full measure. Define $`E_y:=J_{x,y}(F_xE_{x,y})`$. The set $`E_y`$ is of $`W_y`$-measure $`1/2`$, because $`F_xE_{x,y}`$ is of $`W_x`$-measure $`1/2`$ and that $`J_{x,y}`$ sends $`W_x`$ on $`W_y`$. Let $`\gamma =J_{x,y}(\gamma ^{})E_y`$. Using Lemma 4.5, we have
$$\underset{t\mathrm{}}{lim\; sup}d(\gamma (t),\gamma ^{}(t))\epsilon .$$
Thus one gets
$$\underset{n\mathrm{}}{lim\; sup}|B_n(f,\gamma )B_n(f,\gamma ^{})|\alpha ,$$
and the lemma follows because the Birkhoff sums of $`\gamma ^{}`$ converge to $`f𝑑\mu `$. $`\mathrm{}`$
We are now able to finish the proof of the theorem. Let $`\mu ^{}`$ be another ergodic measure, and $`f`$ be a continuous function on $`M`$. We are going to prove that $`f𝑑\mu ^{}=f𝑑\mu `$. Observe that by ergodicity of $`\mu ^{}`$ there exists a point $`y`$ on $`M`$ such that for $`w_y`$-almost every $`\gamma \mathrm{\Gamma }_y`$, the Birkhoff sums $`B_n(f,\gamma )`$ converge to $`f𝑑\mu ^{}`$. Denote by $`F_y`$ the set of such paths $`\gamma `$ which is of full $`w_y`$-measure. We apply Lemma 4.7 to the point $`y`$: for every $`\gamma E_y`$, the Birkhoff sums $`B_n(f,\gamma )`$ are tending to $`f𝑑\mu `$ with an error of $`\alpha `$. Because the measure of $`E_y`$ is positive, it intersects $`F_y`$. Thus one gets $`|f𝑑\mu ^{}f𝑑\mu |\alpha `$. Because $`\alpha `$ is arbitrary, there is only one ergodic harmonic measure, and thus only one harmonic measure. The theorem is proved. $`\mathrm{}`$
###### Proposition 4.8.
Let $`(,\mathrm{\Delta }_g)`$ be a similar Laplace foliation of codimension $`1`$ of a compact manifold $`M`$, where $`\mathrm{\Delta }_g`$ is the Laplacian of a riemannian metric. Then every ergodic harmonic measure is supported on a minimal subset.
Proof. Let $`\mu `$ be an ergodic harmonic measure on $`(,\mathrm{\Delta }_g)`$, and $``$ be a minimal closed subset contained in the support of $`\mu `$. We are going to prove that $``$ is exactly the support of $`\mu `$. It is clear that one can suppose that $`\mu `$ does not charge any leaf.
Suppose that the foliation is oriented. Because the operator $`\mathrm{\Delta }_g`$ is symmetric, the measure $`\mu `$ induces a transverse harmonic measure which is $`\mathrm{\Delta }_g`$-harmonic. From every point $`x`$ of $`M`$, the positive $`𝒢`$-orbit of $`x`$ intersects $``$ in a first time in a point $`y`$. Consider the function $`f(x)=\nu ([x,y))`$. This is a continuous function, because $`\mu `$ does not charge any leaf, and it is harmonic on every leaf. By Garnett lemma \[Ga\], this function has to be constant on $`\mu `$-almost every leaf, thus by continuity of $`f`$, $`f`$ is constant on the support of $`\mu `$. But on the minimal $``$, $`f`$ vanishes, so that the restriction of $`f`$ to the support of $`\mu `$ is identically $`0`$. Thus the support of $`\mu `$ has to be reduced to $``$. The proposition is proved. $`\mathrm{}`$
###### Example 4.9.
Theorem 4.6 seems to be false when the drift does not preserve the volume $`\mathrm{vol}_g`$. We give an example of a similar Laplace lamination of a compact space which is minimal, transversely conformal, and not uniquely ergodic.
A lamination of a compact space $`X`$ is an atlas of homeomorphisms from open sets of $`X`$ to the product of an euclidian ball by a topological set, in such a way that the changes of coordinates preserve the local fibration by balls, and the diffeomorphisms from a piece of ball to another depends continuously of the transverse parameter in the smooth topology. The definition of a Laplace operator on a lamination is exactly the same as in the foliation case.
We are going to describe an example of a similar Laplace lamination of a compact space, which has been constructed by Sullivan \[Su3\]. Let $`𝐇`$ be the upper-half plane, whose hyperbolic metric is expressed by
$$g=\frac{dx^2+dy^2}{y^2},$$
in the coordinates $`z=x+iy`$ of $`𝐇`$. Consider the unit vector fields that points on the direction $`\mathrm{}`$ of $`𝐇`$. In the $`x,y`$ coordinates, it is expressed as
$$V=y\frac{}{y}.$$
The direct isometries of $`(𝐇,g,V)`$ are the maps of the form $`zaz+b`$ where $`a`$ is a positive number, and $`b`$ is a real number. For any real number $`\kappa `$, these transformations preserve the Laplace operator
$$\mathrm{\Delta }_\kappa =\mathrm{\Delta }_g+\kappa V.$$
Let $`A(𝐙[1/2])`$ be the group of affine transformations $`xax+b`$, where $`a`$ is a power of $`2`$, and $`b`$ is a dyadic integer of the form $`p/2^n`$, $`p`$ and $`n`$ being integers. The group $`A(𝐙[1/2])`$ acts naturally on the product $`𝐇\times 𝐐_2`$ of the upper half plane by the field of $`2`$-adic numbers, the action preserving the natural structure of the horizontal Laplace lamination $`(𝐇\times 𝐐_2,\mathrm{\Delta }_\kappa )`$. The action is discrete, without fixed point, and the quotient $`(X,\mathrm{\Delta }_\kappa )`$ is a similar Laplace lamination of a compact space.
When $`\kappa >1`$ the Laplace laminations $`(X,\mathrm{\Delta }_\kappa )`$ carry harmonic measures that charge any cylinder leaves. This can be seen by the same arguments as those given in 3.4. These examples are not codimension $`1`$ foliations, but they share with codimension $`1`$ foliation the property of being transversely conformal, which is the only property used in the proof of Theorem 4.6.
It seems to us that the hypothesis $`div_{\mathrm{vol}_g}V=0`$ is too strong. For instance, we conjecture that for any Laplace operator on the base of a suspension, the conclusion of 4.6 holds.
There are analog examples of similar Laplace laminations of a compact space, associated with tilings of the hyperbolic plane. They have been studied by Petite \[Pe\]; many of them are minimal but not uniquely ergodic even for a symmetric operator. The lack of unique ergodicity is due to the fact that they do not carry a transversely invariant “conformal” structure.
## 5. Appendix: Technical proofs
### 5.1.
Proof of Proposition 2.3. It is a well-known fact (see \[C-L-Y, Ma\]), that for a Brownian motion, the probability of making large steps decreases very fast:
(5.1)
$$C_1,d_0:pM,d>d_0W_p(\mathrm{dist}(\stackrel{~}{\gamma }(\delta ),p)>d)e^{C_1d}.$$
Thus, we may choose a random variable $`\xi `$, such that $`\xi 0`$, $`E\xi <\mathrm{}`$ and
(5.2)
$$pM,d>0W_x(\mathrm{dist}(\stackrel{~}{\gamma }(\delta ),p)>d)P(\xi >d).$$
Let us choose for every $`x`$, a function $`\chi _x:\mathrm{\Gamma }_x`$, such that $`\chi _x`$ depends only on $`\gamma |_{[0,\delta ]}`$, has the same distribution as $`\xi `$ and such that $`\chi _x(\gamma )k_1(\gamma )`$ for every $`\gamma \mathrm{\Gamma }_x`$. Recall, that $`\sigma :\mathrm{\Gamma }\mathrm{\Gamma }`$ stays for the map, erasing the first step of the $`\delta `$-discretization: $`\sigma (\gamma )(t)=\gamma (t+\delta )`$. Then,
$$j,\gamma k_j(\gamma )=k_1(\sigma ^{j1}(\gamma )).$$
Let us denote $`\xi _j(\gamma )=\chi _{\gamma ((j1)\delta )}(\sigma ^{j1}(\gamma ))`$. Then,
$$\begin{array}{c}\frac{k_1(\gamma )+\mathrm{}+k_n(\gamma )}{n}=\frac{k_1(\gamma )+k_1(\sigma (\gamma ))+\mathrm{}+k_1(\sigma ^{n1}\gamma )}{n}\hfill \\ \hfill \frac{\chi _{\gamma (0)}(\gamma )+\chi _{\gamma (\delta )}(\sigma (\gamma ))+\mathrm{}+\chi _{\gamma ((n1)\delta )}(\sigma ^{n1}\gamma )}{n}=\\ \hfill =\frac{\xi _1(\gamma )+\mathrm{}+\xi _n(\gamma )}{n}\end{array}$$
Note, that all the $`\xi _j`$ are distributed identically with $`\xi `$. Moreover, from the Markovian property, the conditional distribution of the variable $`\xi _{n+1}`$ with respect to every condition $`\gamma |_{[0,n\delta ]}=\overline{\gamma }`$ coincides with the distribution of $`\xi `$. On the contrary, the variables $`\xi _1,\mathrm{},\xi _n`$ are determined by such a condition. Thus, the variable $`\xi _{n+1}`$ is independent from $`\xi _1,\mathrm{},\xi _n`$. As $`n`$ is arbitrary, all the variables $`\xi _1,\mathrm{},\xi _n,\mathrm{}`$ are independent and identically distributed.
Now, let us take $`c=E\xi +1`$. We are going to show that for almost every $`\gamma \mathrm{\Gamma }`$ we have $`\underset{n\mathrm{}}{lim\; sup}K_n(\gamma )/n<c`$. From the Large Numbers Law, we have:
$$\underset{n\mathrm{}}{lim\; sup}\frac{\xi _1(\gamma )+\mathrm{}+\xi _n(\gamma )}{n}=E\xi <c$$
$`W_x`$–almost surely. But
$$\underset{n\mathrm{}}{lim\; sup}\frac{k_1(\gamma )+\mathrm{}+k_n(\gamma )}{n}\underset{n\mathrm{}}{lim\; sup}\frac{\xi _1(\gamma )+\mathrm{}+\xi _n(\gamma )}{n},$$
and thus
$$\underset{n\mathrm{}}{lim\; sup}\frac{k_1(\gamma )+\mathrm{}+k_n(\gamma )}{n}<c$$
$`W_x`$-almost surely. This concludes the proof of the proposition. $`\mathrm{}`$
### 5.2. Proof of Lemma 2.4
First, we are going to prove the lemma in the particular case of a codimension one foliation $``$, using the transversal one-dimensional foliation $`𝒢`$, described in Section 2.
In order to prove the lemma, we shall study the behaviour of heat kernels for a time $`t=\delta `$ fixed on different but close enough leaves. We are going to use the following idea: major parts of the heat distribution measures on these two leaves are similar (i.e. the density of $`𝒢`$-holonomy image of one with respect to another is close to $`1`$). Moreover, the infinite product of these densities converges along most Brownian paths, because along these paths the leaves approach exponentially. Thus, the measures $`W_x^\delta =F_{}W_x|_{E_x}`$ and $`(F\mathrm{\Phi }_{\overline{x},x})_{}W_{\overline{x}}|_{E_{\overline{x}}}`$ are absolutely continuous with respect to each other, and the density on the major part of trajectories is close to $`1`$. In particular, the total measures $`W_x(E_x)`$ and $`W_{\overline{x}}(E_{\overline{x}})`$ are close to each other.
###### Proposition 5.1.
Let $`\overline{x}𝒢_x`$, $`\mathrm{dist}_𝒢(x,\overline{x})=\theta `$. Consider the measures $`\nu _x`$ and $`\nu _{\overline{x}}`$, where the measure $`\nu _z=p(z,;\delta )d\mathrm{vol}_g`$ on $`_z`$ gives the distribution of Brownian motion at the time $`\delta `$. Also, let $`\epsilon _2>0`$ and $`R`$ be chosen. Then, there exists a set $`S=S(R,\epsilon _2,\theta ,x)B_R^{}(x)`$, such that
* $`\nu _x(S)>1\epsilon _1`$,
* $`\frac{d\nu _x}{d(\mathrm{\Phi }_{x,\overline{x}}^{}\nu _{\overline{x}})}|_S[1\epsilon _2,1+\epsilon _2]`$,
where $`\epsilon _1=\frac{1}{\epsilon _2}\left(G_1^{}e^{G_2^{}R}\theta +e^{G_3R^2}\right)`$, and $`G_1^{},G_2^{},G_3`$ are geometric constants (depending only on the foliation $``$).
Proof. First, let us equip the leaf $`_x`$ with another Riemannian metric $`g^{}`$, coinciding with $`\mathrm{\Phi }_{}g|_{_{\overline{x}}}`$ inside $`B_R^{}(x)`$ and with $`g|__x`$ outside $`B_{2R}^{}(x)`$; in the annulus left, the metrics $`g^{}`$ is defined using cut-off function.
Note, that the distance (in $`C^{2d}`$-topology, where $`d`$ is the dimension of the leaves of $``$) between $`g`$ and $`g^{}`$ is at most $`G_0\theta e^{G_12R},`$ where $`G_1`$ is a constant, giving the maximum deviation of leaves of $``$, and $`G_0`$ gives the maximum of derivative of $`g`$ along $`𝒢`$.
Also, we recall, that the heat kernel can be constructed explicitly as a series of convolutions. This procedure is described in the books of Candel \[C-C\] and Chavel \[Cha\]. In a few words, a function $`L`$, which “almost satisfies” the heat equation, is explicitly constructed, and then the real heat kernel is obtained as a sum of $`L`$ and a series of convolutions. These series are converging uniformly for every fixed moment of time $`t`$, and the dependence of the metrics is smooth (due to the explicit nature of the construction). Thus, the distance between the heat kernels for the metrics $`g`$ and $`g^{}`$ at the time $`t=\delta `$ can be bounded by the product of a constant $`G_2`$ (depending only on the geometry of foliation $``$ and of the moment $`\delta `$) and of the distance between $`g`$ and $`g^{}`$. This distance is at most $`e^{2G_1R}\theta `$; hence, the difference between these kernels is bounded by $`G_2e^{2G_1R}\theta .`$
Now, due to the upper bounds for the heat kernel \[C-L-Y, Ma\], the set of Brownian trajectories on the interval of time $`[0,\delta ]`$, starting at $`x`$ and exiting from the ball $`B_R^{}(x)`$ at some intermediate moment, has the measure at most $`e^{G_3R^2}`$. We notice that these trajectories for the metric $`g^{}`$ and for the metric $`\mathrm{\Phi }_{x,\overline{x}}^{}g|_{_{\overline{x}}}`$ are the same (they do not pass through the points where these metrics do not coincide). Thus, the parts of heat kernel at time $`\delta `$, coming from these trajectories, are the same for these two metrics.
Let us define the set $`S`$ in a following way: $`zS`$, if
* The density $`p_g^{}(x,z;\delta )`$ is at least $`3G_2e^{2G_1R}\theta /\epsilon _2`$
* At least $`1\epsilon _2/3`$ of this density comes from the trajectories staying inside $`B_R^{}`$ (in particular, $`zB_R^{}(x)`$).
Then, the second condition for $`S`$ (quotient of densities) is satisfied automatically: the maximum possible change of density at a point of $`S`$ is the sum of changes while passing from $`g`$ to $`g^{}`$ (at most $`\epsilon _2/3`$ part) and from $`g^{}`$ to $`g|_{_{\overline{x}}}`$ (at most $`\epsilon _2/3`$ part due to common set of trajectories).
Now, let us estimate $`\nu (_xS)`$, thus verifying the first condition. The points of this complementary can be of two types: either points of $`B_R^{}`$ with too small value of density (we note this set $`X_1`$), or with too big part of this density coming from trajectories, exiting the ball $`B_R^{}(x)`$ (we note this set $`X_2`$).
The first part is estimated as
$$\begin{array}{c}\nu (X_1)=_{B_R^{}(x)}p(x,z;\delta )𝑑\mathrm{vol}_g3G_2e^{2G_1R}\theta /\epsilon _2\mathrm{vol}_g(B_R^{}(x))\hfill \\ \hfill (3G_2e^{2G_1R}\theta /\epsilon _2)e^{G_4R},\end{array}$$
where $`G_4`$ is the constant, bounding the growth of the leaves of $``$.
The second part is estimated as follows: denote by $`\rho (z)`$ the part of the density $`p(x,z;\delta )`$, coming from the trajectories exiting from the ball $`B_R^{}(x)`$. Then,
$$\begin{array}{c}\nu (X_2)=_{\{z:\rho (z)/p(x,z;\delta )>\epsilon _2/3\}}p(x,z;\delta )𝑑\mathrm{vol}_g(z)\hfill \\ \hfill _{\{z:\rho (z)/p(x,z;\delta )>\epsilon _2/3\}}\frac{3}{\epsilon _2}\rho (z)𝑑\mathrm{vol}_g(z)\\ \hfill \frac{3}{\epsilon _2}__x\rho (z)𝑑\mathrm{vol}_g(z)<\frac{3}{\epsilon _2}e^{G_3R^2}\end{array}$$
(the last inequality comes from the upper bound for the probability of all the set of trajectories, leaving the ball of radius $`R`$ at some moment between 0 and $`\delta `$).
Finally, we obtain
$$\begin{array}{c}\nu (_xS)=\nu (X_1X_2)<\frac{3G_2e^{(2G_1+G_4)R}}{\epsilon _2}\theta +\frac{3}{\epsilon _2}e^{G_3R^2}=\hfill \\ \hfill =\frac{1}{\epsilon _2}\left(G_1^{}e^{G_2^{}R}\theta +e^{G_3R^2}\right),\end{array}$$
where $`G_1^{}=3G_2`$, $`G_2^{}=2G_1+G_4`$.
The first condition on $`S`$ is satisfied. $`\mathrm{}`$
For any $`R,\theta >0`$ let us denote
$$\mathrm{\Psi }(R,\theta )=\sqrt{G_1^{}e^{G_2^{}R}\theta +e^{G_3R^2}}.$$
Also, let $`r(\theta )=(\mathrm{log}\frac{1}{2\theta G_1^{}})/(2G_2^{})`$. Then $`G_1^{}e^{G_2^{}r(\theta )}\theta =\frac{1}{2}\sqrt{\theta }`$, and for all $`\theta `$ sufficiently small $`e^{G_3r(\theta )^2}<\frac{1}{2}\sqrt{\theta }`$. Thus, for all $`\theta `$ sufficiently small
$$\mathrm{\Psi }(\theta ):=\mathrm{\Psi }(r(\theta ),\theta )<\sqrt{\theta }.$$
Now, denote
$$S(\theta ,x)=S(r(\theta ),\mathrm{\Psi }(\theta ),\theta ,x).$$
For this set, the conclusions of Proposition 5.1 is satisfied with $`\epsilon _1=\epsilon _2<\sqrt{\theta }`$. Let us choose a small transversal interval $`JI`$, $`Jx`$, and consider the subset of $`E_x`$, defined as
(5.3)
$$E_x^{}=\left\{\gamma E_x\right|n0x_{n+1}S(\theta _n,x_n)\},$$
where $`x_n=\gamma (n\delta )`$ is the discretization sequence, corresponding to $`\gamma `$, and $`\theta _n=|h_{\gamma |_{[0,n\delta ]}}(J)|`$ is the sequence of the corresponding transverse distances (exponentially decreasing due to the nature of $`E_x`$).
###### Lemma 5.2.
For $`\overline{x}J`$, the images of the measure $`W_x|_{E_x^{}}`$ under $`F`$ and of the measure $`W_{\overline{x}}|_{E_{\overline{x}}^{}}`$ under $`F\mathrm{\Phi }_{\overline{x},x}`$ are absolutely continuous with respect to each other. The density can be made arbitrary close to 1 by choice of sufficiently small intervals $`I`$ and $`J`$. Moreover, by such a choice, the differences $`W_x(E_x)W_x(E_x^{})`$ and $`W_{\overline{x}}(E_{\overline{x}})W_{\overline{x}}(E_{\overline{x}}^{})`$ can be made arbitrarily small.
Proof. Let us consider the projection map
$$\pi _n:(\stackrel{~}{}_x)^{\mathrm{}}(\stackrel{~}{}_x)^{n+1},\pi _n(\{x_j\}_{j=0}^{\mathrm{}})=\{x_j\}_{j=0}^n$$
and its composition with $`F`$, which we denote $`F_n:E_x^{}(\stackrel{~}{}_x)^{n+1}`$,
$$F_n(\gamma )=\{\stackrel{~}{\gamma }(j\delta )\}_{j=0}^n.$$
Let us denote
$$\mu _1=F_{}W_x|_{E_x^{}},\mu _2=(F\mathrm{\Phi }_{\overline{x},x})_{}W_{\overline{x}}|_{E_{\overline{x}}^{}},$$
$$\mu _1^n=(\pi _n)_{}\mu _1=(F_n)_{}W_x|_{E_x^{}},\mu _2^n=(\pi _n)_{}\mu _2=(F_n\mathrm{\Phi }_{\overline{x},x})_{}W_{\overline{x}}|_{E_{\overline{x}}^{}}.$$
###### Proposition 5.3.
The measures $`\mu _1^n`$ and $`\mu _2^n`$ are absolutely continuous with respect to each other, and
$$\frac{d\mu _2^n}{d\mu _1^n}|_{(x_j)}=\rho _n((x_j))=\underset{j=0}{\overset{n1}{}}\frac{d\nu _{x_j}}{d(\mathrm{\Phi }_{x_j,\overline{x}_j}^{}\nu _{\overline{x}_j})},$$
To prove Lemma 5.2, it suffices to show that for $`\mu _1`$–almost every point $`(x_j)(\stackrel{~}{}_x)^{}`$ the sequence $`\rho _n((x_j))`$ converges to some number between $`0`$ and $`\mathrm{}`$. It is equivalent to the convergence of the infinite product
(5.4)
$$\underset{j=1}{\overset{\mathrm{}}{}}\frac{d\nu _{x_j}}{d(\mathrm{\Phi }_{x_j,\overline{x}_j}^{}\nu _{\overline{x}_j})}=\underset{n\mathrm{}}{lim}\rho _n((x_j))$$
(again, for $`\mu _1`$–almost every $`(x_j)(\stackrel{~}{}_x)^{}`$).
Now, for every $`(x_j)`$ in the image $`F(E_x^{})`$, the logarithm of the product (5.4) can be estimated as
(5.5)
$$\begin{array}{c}\left|\mathrm{log}\underset{j=1}{\overset{\mathrm{}}{}}\frac{d\nu _{x_j}}{d(\mathrm{\Phi }_{x_j,\overline{x}_j}^{}\nu _{\overline{x}_j})}\right|\underset{j=0}{\overset{\mathrm{}}{}}\left|\mathrm{log}\frac{d\nu _{x_j}}{d(\mathrm{\Phi }_{x_j,\overline{x}_j}^{}\nu _{\overline{x}_j})}\right|\hfill \\ \hfill \underset{j=0}{\overset{\mathrm{}}{}}\left|2\sqrt{\theta _j}\right|\underset{j=0}{\overset{\mathrm{}}{}}\left|2\sqrt{Ce^{\alpha j}|J|}\right|C_3\sqrt{|J|}\underset{j=0}{\overset{\mathrm{}}{}}e^{\alpha j/2}=C_4\sqrt{|J|}.\end{array}$$
Here we used the definition of $`E_x^{}`$ to bound the density $`\frac{d\nu _{x_j}}{d(\mathrm{\Phi }_{x_j,\overline{x}_j}^{}\nu _{\overline{x}_j})}`$, and then again to estimate $`\theta _n`$.
We have estimated the density; moreover, for $`J`$ sufficiently small this density (due to (5.5)) can be made arbitrary close to $`1`$.
Now, let us estimate the difference $`W_z(E_z)W_z(E_z^{})`$, where $`zJ`$. Denote
$$E_{z,n}=\left\{\gamma E_z\right|j,0j<nx_{n+1}S(\theta _n,x_n)\},$$
$$\stackrel{~}{C}_{z,n}=F_n(E_{z,n}),$$
and let $`\stackrel{~}{\mu }_n`$ be a measure on $`(\stackrel{~}{}_z)^{n+1}`$, defined as the discretization image of $`W_z|_{E_{z,n}}`$. Also, consider projection maps $`\stackrel{~}{\pi }_n:\stackrel{~}{C}_{z,n+1}\stackrel{~}{C}_{z,n}`$. Then, $`(\stackrel{~}{\pi }_n)_{}\stackrel{~}{\mu }_{n+1}`$ is absolutely continuous with respect to $`\stackrel{~}{\mu }_n`$, and the density is equal to
$$\stackrel{~}{\rho }_n((x_j))=\nu _{x_n}(S(\theta _n,x_n))12\sqrt{\theta _n}12\sqrt{Ce^{\alpha n}|J|}.$$
Thus,
$$\begin{array}{c}W_z(E_{n,z}E_{n+1,z})=_{\stackrel{~}{C}_{n,z}}(1\stackrel{~}{\rho }_n)((x_j))𝑑\stackrel{~}{\mu }_n((x_j)_{j=0}^n)\hfill \\ \hfill \underset{\stackrel{~}{C}_{n,z}}{}2\sqrt{Ce^{\alpha n}|J|}𝑑\stackrel{~}{\mu }_n((x_j)_{j=0}^n)C_3e^{\alpha n/2}\sqrt{|J|}.\end{array}$$
Now, recall that $`E_z^{}=_{n=1}^{\mathrm{}}E_{z,n}`$, hence,
(5.6)
$$\begin{array}{c}W_z(E_z)W_z(E_z^{})=\underset{n=0}{\overset{\mathrm{}}{}}W_z(E_{z,n}E_{z,n+1})\hfill \\ \hfill \underset{n=0}{\overset{\mathrm{}}{}}C_3e^{\alpha n/2}\sqrt{|J|}C_4\sqrt{|J|}.\end{array}$$
The difference $`W_z(E_z)W_z(E_z^{})`$ tends to 0 as $`|J|`$ tends to $`0`$. $`\mathrm{}`$
Proof of Lemma 2.4.
#### 5.2.1. “Positive measure” part
First, let us prove the “positive measure” part. Namely, estimate the difference $`|W_x(E_x)W_{\overline{x}}(E_{\overline{x}})|`$ if $`\overline{x}JI`$:
$$\begin{array}{c}|W_x(E_x)W_{\overline{x}}(E_{\overline{x}})||W_x(E_x)W_x(E_x^{})|+\hfill \\ \hfill +|W_x(E_x^{})W_{\overline{x}}(E_{\overline{x}}^{})|+|W_{\overline{x}}(E_{\overline{x}}^{})W_{\overline{x}}(E_{\overline{x}})|.\end{array}$$
All the three differences can be estimated using Lemma 5.2. Thus, choosing any $`p_1<W_x(E_x)`$, we can find a sufficiently small transversal interval $`J`$, such that for any $`\overline{x}J`$ we have $`W_{\overline{x}}(E_{\overline{x}})p_1`$.
Now, suppose that $`y\stackrel{~}{}_{\overline{x}}`$. Then, the measures $`\nu _y`$ and $`\nu _{\overline{x}}`$ are absolutely continuous with respect to each other, and the density on the major (with respect to these measures) part of $`\stackrel{~}{}_{\overline{x}}`$ is close to 1. Let us choose $`\epsilon _3>0`$ and denote the set
$$\stackrel{~}{S}(\overline{x},y,\epsilon _3)=\{z\stackrel{~}{}_{\overline{x}}\frac{d\nu _{\overline{x}}}{d\nu _y}(z)[\frac{1}{1+\epsilon _3},1+\epsilon _3]\}.$$
Then,
$$\epsilon _3>0r>0:y\stackrel{~}{}_{\overline{x}},d(y,\overline{x})<r\nu _{\overline{x}}(\stackrel{~}{S})>1\epsilon _3,\nu _y(\stackrel{~}{S})>1\epsilon _3.$$
Note that due to the Markovian property (the conditional distribution of $`\sigma (\gamma )`$ with respect to every condition $`\sigma (\gamma )(0)=z^{}`$ coincides with $`W_z^{}`$), we have
$$W_z(E_z)=_{\stackrel{~}{}_z}W_z^{}(E_z^{}^1)𝑑\nu _z(z^{}),$$
where
$$E_z^{}^1=\{\gamma n\left|h_{\gamma |_{[0,n\delta ]}}(h_{z,z^{}}(I))\right|<C_0e^{\alpha (n+1)\delta }|I|\}$$
(this is the definition of $`E_z`$, rewritten in terms of the shift $`\sigma (\gamma )`$). Thus, for $`d(y,\overline{x})<r`$ we have
(5.7)
$$W_y(E_y)=W_y(E_y\{x_1\stackrel{~}{S}\})+_{\stackrel{~}{S}}W_z(E_z^1)𝑑\nu _y(z),$$
(5.8)
$$\begin{array}{c}W_{\overline{x}}(E_{\overline{x}})=W_{\overline{x}}(E_{\overline{x}}\{x_1\stackrel{~}{S}\})+_{\stackrel{~}{S}}W_z(E_z^1)𝑑\nu _{\overline{x}}(z)=\hfill \\ \hfill =W_{\overline{x}}(E_{\overline{x}}\{x_1\stackrel{~}{S}\})+_{\stackrel{~}{S}}W_z(E_z^1)\frac{d\nu _{\overline{x}}}{d\nu _y}(z)𝑑\nu _y(z).\end{array}$$
The first summands in the right hand sides of (5.7) and (5.8) are no greater than $`\epsilon _3`$, and the quotient of the second summands is bounded by $`1+\epsilon _3`$. Thus, for $`y`$ and $`\overline{x}`$ being sufficiently close to each other, the probabilities $`W_{\overline{x}}(E_{\overline{x}})`$ and $`W_y(E_y)`$ are also close to each other. In particular, for every $`p_0<p_1`$ we can find $`r>0`$, such that for $`U`$ being the union of $`r`$-leafwise-neighborhoods of points of $`J`$, we have
$$yUW_y(E_y)p_0.$$
It completes the proof of this part of the lemma.
#### 5.2.2. Proof of the “distributions” part
Note that due to the arguments already used we may suppose that $`y𝒢_x`$: for $`y`$ and $`y^{}`$ on the same $``$-leaf, the measures $`\nu _y^{}`$ and $`\nu _y`$ are absolutely continuous with respect to each other, and hence any tail-type property holds (or does not hold) simultaneously for typical trajectories in $`\mathrm{\Gamma }_y`$ and $`\mathrm{\Gamma }_y^{}`$. Hence, if $`y`$ does not belong to $`𝒢_x`$, we may replace it by $`y^{}`$, which is an intersection point of $`_y`$ and $`𝒢_x`$.
Recall that almost every trajectory $`\gamma \mathrm{\Gamma }_x`$ (due to the choice of $`x`$) is distributed with respect to $`\mu `$. Thus, for almost every path $`\gamma \mathrm{\Gamma }_x`$ we have
(5.9)
$$\underset{T\mathrm{}}{lim}\frac{1}{T}\gamma _{}\mathrm{leb}_{[0,T]}=\mu .$$
We know that a trajectory of $`E_y`$ approaches a trajectory of $`E_x`$. Unfortunately, we can not claim that the map giving the trajectory of $`E_x`$ by a trajectory of $`E_y`$ is absolutely continuous (or, what is the same, maps typical trajectories to typical ones): we have this statement only for discretizations of trajectories.
Thus, we have to prove the distributions property using discretizations behaviour. The following arguments are a technical realization of this idea.
Rewrite (5.9) in the terms of discretization. Let a continuous function $`\phi `$ on $`M`$ be chosen. Then, for $`W_x`$-almost every $`\gamma \mathrm{\Gamma }_x`$,
(5.10)
$$\left(\frac{1}{T}\gamma _{}\mathrm{leb}_{[0,T]}\right)(\phi )=\frac{1}{T}_0^T\phi (\gamma (t))𝑑t.$$
It is clear that we can restrict to the moments of time of the form $`T=n\delta `$; for such $`T`$,
(5.11)
$$\begin{array}{c}\frac{1}{n\delta }_0^{n\delta }\phi (\gamma (t))𝑑t=\frac{1}{n\delta }\underset{j=0}{\overset{n1}{}}_{j\delta }^{(j+1)\delta }\phi (\gamma (t))𝑑t=\hfill \\ \hfill =\frac{1}{n}\underset{j=0}{\overset{n1}{}}\frac{1}{\delta }_{j\delta }^{(j+1)\delta }\phi (\gamma (t))𝑑t.\end{array}$$
Let $`zM`$ be some point, and let us rewrite (5.11) for a typical trajectory $`\gamma \mathrm{\Gamma }_z`$. Namely, we divide this sum into discrete averaging and the rest term:
(5.12)
$$\begin{array}{c}\frac{1}{n\delta }_0^{n\delta }\phi (\gamma (t))𝑑t=\frac{1}{n}\underset{j=0}{\overset{n1}{}}\phi (\gamma (j\delta ))+\hfill \\ \hfill +\frac{1}{n}\underset{j=0}{\overset{n1}{}}\frac{1}{\delta }_{j\delta }^{(j+1)\delta }(\phi (\gamma (t))\phi (\gamma (j\delta )))𝑑t.\end{array}$$
We estimate the second term in the right hand side of (5.11), decomposing the sum in two parts, the one corresponding to $`j`$ with $`diam(\gamma ([j\delta ,(j+1)\delta ]))<r`$ and the one with $`diam(\gamma ([j\delta ,(j+1)\delta ]))r`$.
(5.13)
$$\begin{array}{c}|\frac{1}{n}\underset{j=0}{\overset{n1}{}}\frac{1}{\delta }_{j\delta }^{(j+1)\delta }(\phi (\gamma (t))\phi (\gamma (j\delta )))dt.|\hfill \\ \hfill \frac{1}{n}\underset{j<n,diam(\gamma ([j\delta ,(j+1)\delta ]))<r}{}\frac{1}{\delta }_{j\delta ^{\prime \prime }}^{(j+1)\delta ^{\prime \prime }}\left|\phi (\gamma (j\delta ))\phi (\gamma (t))dt\right|+\\ \hfill +\frac{1}{n}\underset{j<n,diam(\gamma ([j\delta ,(j+1)\delta ]))r}{}\frac{1}{\delta }_{j\delta ^{\prime \prime }}^{(j+1)\delta ^{\prime \prime }}\left|\phi (\gamma (j\delta ))\phi (\gamma (t))dt\right|\\ \hfill \omega _\phi (r)+2\frac{\mathrm{\#}\{j:diamr\}}{n}\underset{}{sup}|\phi |,\end{array}$$
where $`\omega _\phi `$ is the modulus of continuity of the function $`\phi `$.
Let us pass in (5.13) to the upper limit:
(5.14)
$$\begin{array}{c}\underset{n\mathrm{}}{lim\; sup}|\frac{1}{n}\underset{j=0}{\overset{n1}{}}\frac{1}{\delta }_{n\delta }^{(j+1)\delta }(\phi (\gamma (t))\phi (\gamma (j\delta )))dt.|\hfill \\ \hfill \omega _\phi (r)+2\underset{}{sup}|\phi |\underset{j\mathrm{}}{lim\; sup}\frac{\mathrm{\#}\{j<ndiam(\gamma [j\delta ,(j+1)\delta ])r\}}{n}.\end{array}$$
For a function $`\phi `$ chosen, the first summand can be made arbitrarily small by a choice of $`r`$ due to the continuity of $`\phi `$. For every $`r>0`$ chosen, the second summand can be made arbitrarily small for almost all trajectories by a choice of sufficiently small $`\delta >0`$ (which is uniform in $`zM`$) due to the same arguments as the ones used in the proof of Proposition 2.3.
Hence, for every $`\epsilon >0`$ we can find $`r`$ and then $`\delta `$ small enough, such that for almost all $`\gamma \mathrm{\Gamma }_z`$,
$$\underset{n\mathrm{}}{lim\; sup}\left|\frac{1}{n}\underset{j=0}{\overset{n1}{}}\frac{1}{\delta }_{j\delta }^{(j+1)\delta }(\phi (\gamma (t))\phi (\gamma (j\delta )))𝑑t\right|<\epsilon .$$
The rest term in (5.12) is estimated, and taking it together with (5.10), for a discretization $`(x_j)=F^\delta (\gamma )`$ of a typical path $`\gamma \mathrm{\Gamma }_x`$, we have
$$\underset{n\mathrm{}}{lim\; sup}\left|\frac{1}{n}\underset{j=0}{\overset{n1}{}}\phi (x_j)_M\phi 𝑑\mu \right|<\epsilon .$$
Now, for a $`W_{\overline{x}}`$-typical path $`\gamma `$ from $`E_y^{}`$ let us estimate the difference
(5.15)
$$\underset{n\mathrm{}}{lim\; sup}\left|\frac{1}{n\delta }_0^{n\delta }\phi (\gamma (t))𝑑t_{}\phi 𝑑\mu \right|.$$
We have:
$$\begin{array}{c}\underset{n\mathrm{}}{lim\; sup}\left|\frac{1}{n\delta }_0^{n\delta }\phi (\gamma (t))𝑑t_{}\phi 𝑑\mu \right|\hfill \\ \hfill \underset{n\mathrm{}}{lim\; sup}\left|\frac{1}{n\delta }_0^{n\delta }\phi (\gamma (t))𝑑t\frac{1}{n}\underset{j=0}{\overset{n1}{}}\phi (y_j)\right|+\\ \hfill +\underset{n\mathrm{}}{lim\; sup}\frac{1}{n}\underset{j=k}{\overset{n1}{}}\left|\phi (y_j)\phi (x_j)\right|+\\ \hfill +\underset{n\mathrm{}}{lim\; sup}\left|\frac{1}{n}\underset{j=0}{\overset{n1}{}}\phi (x_j)_{}\phi 𝑑\mu \right|,\end{array}$$
where $`(y_n)`$ is the discretization of the path $`\gamma \mathrm{\Gamma }_y`$, and $`(x_n)`$ is its $`\mathrm{\Phi }_{y,x}`$-image. We know that the measures $`F_{}^\delta W_x|_{E_x^{}}`$ and $`(F^\delta \mathrm{\Phi }_{y,x})_{}W_y|_{E_y^{}}`$ are absolutely continuous; thus, the image of a typical sequence is a typical sequence. Hence, for a typical trajectory $`\gamma E_y^{}`$ the first and the last summands do not exceed $`\epsilon `$. So, we have
(5.16)
$$\underset{n\mathrm{}}{lim\; sup}\left|\frac{1}{n\delta }_0^{n\delta }\phi (\gamma (t))𝑑t_{}\phi 𝑑\mu \right|<2\epsilon $$
for a typical path $`\gamma E_y^{}`$, which implies, that
(5.17)
$$\underset{T\mathrm{}}{lim\; sup}\left|\frac{1}{T}_0^T\phi (\gamma (t))𝑑t_{}\phi 𝑑\mu \right|<2\epsilon .$$
For a typical path from $`E_y^{}`$ we have obtained an estimate on the difference between the integral of $`\phi `$ and its average along the path. Let us extend this statement to all the $`E_y`$. Namely, repeating the arguments used in the proof of Remark 2.8, we see that almost every path from $`E_y`$ can be decomposed into a finite starting segment and a path from some $`E_z^{}`$ for some $`z`$ close to $`M`$. Then, the estimate (5.17) holds also for a $`W_y`$-typical path from $`E_y`$.
But the definition of $`E_y`$ does not depend on $`\delta `$, thus, choosing arbitrarily small $`\delta `$ and $`r`$, we have finally:
$$\underset{n\mathrm{}}{lim\; sup}\left|\frac{1}{T}_0^T\phi (\gamma (t))𝑑t_{}\phi 𝑑\mu \right|=0.$$
Hence,
$$\frac{1}{T}_0^T\phi (\gamma (t))𝑑t\underset{T\mathrm{}}{\overset{}{}}\phi 𝑑\mu $$
for a typical path $`\gamma E_y`$.
Recall that by definition the measures $`\mu _t`$ weakly converge to $`\mu `$ if and only if for every continuous function $`\phi `$ we have
$$\phi 𝑑\mu _t\phi 𝑑\mu .$$
Moreover, it suffices to check such convergence for a well-chosen countable family $`\phi _k`$. We have already obtained, that for any function $`\phi `$ and for a typical path $`\gamma E_y`$
$$\left(\frac{1}{T}\gamma _{}\mathrm{leb}_{[0,T]}\right)(\phi )=\frac{1}{T}_0^T\phi (\gamma (t))𝑑t\underset{T\mathrm{}}{\overset{}{}}\phi 𝑑\mu .$$
A countable family of typically satisfied conditions still is a typically satisfied condition, and hence for $`W_y`$-almost every trajectory $`\gamma E_y`$
$$\underset{t\mathrm{}}{lim}\frac{1}{t}\gamma _{}\mathrm{leb}_{[0,t]}=\mu .$$
This completes the proof of the lemma. $`\mathrm{}`$
### 5.3. Codimension higher than one
Here we present a construction which permits us to handle the case of transversely conformal foliation of codimension higher than one. Let $``$ be such a foliation. Equip $`M`$ with a Riemannian metric and for any point $`xM`$, let a transversal $`𝒢_x`$ be an image of the image of the exponential map for a small disk in $`(T_x)^{}`$ disk. We notice that these transversals depend smoothly on $`x`$, and in a small neighborhood of every point $`x`$ for every point $`y`$ there exists a unique point $`z`$ in its small $``$-leafwise neighborhood, such that $`z𝒢_y`$.
We remark that the family of transversals $`\{𝒢_x\}_{xM}`$ does not necessarily form a foliation, for the reason that we can not control intersections of transversals with starting points on different $``$-leaves.
Now, for a point $`x_0`$, let us consider the set
$$\overline{M}=_{y\stackrel{~}{}_{x_0}}\{y\}\times 𝒢_y.$$
Note that $`\overline{M}`$ is a manifold with boundary, naturally inheriting from $`M`$ its foliation structure (except for the fact that for $`\overline{M}`$ some leaves intersect the boundary $`\overline{M}`$) and Riemannian metric on the leaves. But, on $`\overline{M}`$ we have a natural smooth transversal foliation $`𝒢`$, leaves of which are $`\{y\}\times 𝒢_y`$. Now, we can apply the same arguments as these used in the codimension one case, to prove Lemma 2.4.
### 5.4. Non-positivity of Lyapunov exponents
This section is devoted to the following lemma:
###### Lemma 5.4.
Let $``$ be a transversely conformal foliation, $``$ be a minimal set in $``$, and $`\mu `$ be a harmonic ergodic measure supported on $``$. Then, $`\lambda (\mu )0`$.
Proof. We will present the proof for the case of codimension one foliation, using the existence of a transversal foliation $`𝒢`$. Then, it is generalized to an arbitrary codimension case in the same way as in the proof of Lemma 2.4.
Assume the contrary: let $`\lambda (\mu )>0`$. We take some $`\alpha ,\beta `$, $`\alpha <\beta <\lambda (\mu )`$. Prove first that the measure $`\mu `$ does not charge any leaf. Assume the contrary: $`\mu (L)>0`$ for some leaf $`L`$. Then, the density with respect to the volume $`\frac{d\mu |_L}{d\mathrm{vol}_g}`$ is a harmonic function on the leaf $`L`$. Moreover, this function is positive, bounded (because of harmonicity and boundedness of geometry of $`L`$) and of integral $`1`$. Extending this function by 0 to the complementary, we obtain a harmonic measurable positive leafwise integrable function. Garnett \[Ga, Proposition 1, p. 295\] have proved that such a function should be leafwise constant. Thus, it is equal to a positive constant on $`L`$ and hence (as it is integrable) $`L`$ is a compact leaf. But for a compact leaf the Lyapunov exponent equals $`0`$, as the corresponding Dirac measure is transversely invariant.
Let us choose a point $`x_0`$, typical in the sense of Lyapunov exponents: for $`W_{x_0}`$-almost every path $`\gamma \mathrm{\Gamma }_{x_0}`$ the corresponding path has the Lyapunov exponent equal to $`\lambda (\mu )`$.
First, let us consider the simplest case: suppose, that the foliation $`𝒢`$ preserves the metric $`g`$. Also, we suppose that any $``$-along holonomy extends to some fixed neighborhood $`U𝒢_{x_0}`$. Finally, we suppose that the measure $`\mu `$ does not charge any leaf, or equivalently, that measures $`\nu _{}`$ have no points of positive measure.
In this case, the measure $`\mu `$ induces on every leaf $`𝒢_x`$ a conditional measure $`\nu _x`$, which is harmonic in the sense of measure-valued functions (see Section 4.2). Let us take some $`T>0`$ and $`C>1`$ and for every path $`\gamma `$, starting at $`x_0`$, try to find a $`\tau =\tau (\gamma )`$, $`T<\tau <2T`$, as a minimal value $`t_0`$ in this interval possessing the following property:
(5.18)
$$t[0,\tau ]h_{\gamma |_{[\tau ,t]}}^{}(\gamma (\tau ))<Ce^{\beta (\tau t)}.$$
Here, we use $`h_{\gamma |_{[\tau ,t]}}`$ as a short notation for $`h_{\gamma |_{[t,\tau ]}}^1`$ ($`t<\tau `$, and thus the first notation is not absolutely clear). Note, that as the holonomy is taken in the inverse sense, from the moment $`\tau `$ to $`t<\tau `$, one can expect that the derivative is will be small.
Note that the non-existence of such $`\tau `$ means, that
$$\frac{h_{\gamma |_{[0,2T]}}^{}}{h_{\gamma |_{[0,T]}}^{}}e^{\beta T}(x_0),$$
so for $`W_{x_0}`$-almost every path $`\gamma `$ and for every $`T`$, sufficiently big such $`\tau `$ exists. If in the interval $`[T,2T]`$ we can not find a moment satisfying (5.18), then we choose $`\tau (\gamma )=2T`$. Finally, we remark that $`\tau ()`$ is a Markovian moment.
The transversal measures $`\nu _x`$ depend on $`x`$ in a harmonic way, thus for a transversal interval $`I𝒢_{x_0}`$ the measures of its holonomy images $`\nu _{\gamma (t)}(h_{\gamma |_{[0,t]}}(I))`$ form a martingale. The expectation of value of this martingale at the Markovian moment $`\tau `$ should be equal to its initial value:
(5.19)
$$𝔼\nu _{\gamma (\tau (\gamma ))}(h_{\gamma |_{[0,\tau (\gamma )]}}(I))=\nu _{x_0}(I).$$
Now, note that due to the definition of $`\tau (I)`$ for all the paths $`\gamma `$ with $`\tau (\gamma )<2T`$ a neighborhood $`V`$ of $`\gamma (\tau )`$ is contracted exponentially by the holonomy $`h_{\gamma |_{[\tau ,0]}}`$, the radius of $`V`$ is bounded from below by means of $`C`$, $`\alpha ,\beta `$ and the geometry of the foliation. Namely,
$$|h_{\gamma |_{[\tau ,0]}}(U_\epsilon (\gamma (\tau )))|e^{\alpha \tau }e^{\alpha T},$$
where $`\epsilon >0`$ does not depend on $`T`$. Passing from inverse to direct time we see that an exponentially small neighborhood of $`x_0=\gamma (0)`$ is expanded:
$$h_{\gamma |_{[0,\tau ]}}(U_{e^{\alpha T}}(x_0))U_\epsilon (\gamma (\tau )).$$
Now, let us estimate the left part of (5.19) for $`I=U_{e^{\alpha T}}(x_0)`$:
(5.20)
$$\begin{array}{c}𝔼\nu _{\gamma (\tau (\gamma ))}(h_{\gamma |_{[0,\tau (\gamma )]}}(I))\hfill \\ \hfill _{\{\gamma :\tau (\gamma )<2T\}}\nu _{\gamma (\tau (\gamma ))}(h_{\gamma |_{[0,\tau (\gamma )]}}(I))𝑑W_{x_0}\\ \hfill _{\{\gamma :\tau (\gamma )<2T\}}\nu _{\gamma (\tau (\gamma ))}(U_\epsilon (\gamma (\tau )))𝑑W_{x_0}\\ \hfill W_{x_0}\{\gamma :\tau (\gamma )<2T\}\underset{y}{inf}(\nu _y(U_\epsilon (y))).\end{array}$$
As $``$ is compact and $`supp\mu =`$, the infimum in the last term of (5.20) is positive. Thus, the last term stays separated from 0 as $`T`$ tends to infinity, so it is no less than some constant $`c_0>0`$. Hence (5.19) and (5.20) imply that
$$\nu _{x_0}(U_{e^{\alpha T}}(x_0))W_{x_0}\{\gamma :\tau (\gamma )<2T\}\underset{y}{inf}(\nu _y(U_\epsilon (y)))>c_0,$$
and thus the left term does not tend to $`0`$ as $`T`$ tends to infinity. This contradicts the fact that the measure $`\nu _{x_0}`$ can not have atoms. We have obtained the desired contradiction. So, this case is handled.
Let us now consider the case of generic Riemannian structure (not necessarily preserved by the transversal foliation). Note, that the harmonic measure $`\mu `$ still defines conditional measures on the transversals $`\{𝒢_y\}`$, which are its Fubini conditional measures with respect to $`\mathrm{vol}_g`$ on the leaves. We still suppose that the $``$-along holonomy maps are defined on the entire transverse leaves $`\{𝒢_y\}`$. Also, we add the following (simplifying the explanation of this step) hypothesis: all the leaves of $``$ are simply connected.
For these conditional distributions, the harmonicity condition implies that for a transversal interval $`I`$ at a point $`xM`$ and a function $`\rho `$, we have:
(5.21)
$$_I\rho (y)𝑑\nu _x(y)=𝔼_{h_{\gamma |_{[0,t]}}(I)}\frac{p(h_{\gamma |_{[t,0]}}(z),z;t)}{p(x,\gamma (t);t)}\rho (h_{\gamma |_{[t,0]}}(z))𝑑\nu _{\gamma (t)}(z),$$
where the expectation is taken in the sense of $`W_x`$. To prove this formula, we take a smooth function $`f`$ on $`M`$ supported in the neighborhood of $`I`$, with integral on $`_y`$ equal to $`\rho (y)`$. For every fixed $`t`$, as the support of $`f`$ tends to $`I`$, the integral of $`f`$ with respect to $`\mu `$ tends to the left hand side of (5.21), and the integral of $`D^tf`$ to the right hand side. The harmonicity of $`\mu `$ implies that these two integrals coincide, which proves the formula.
Now, let us repeat the arguments used in the similar case with the following modification: we consider only discrete moments of time $`t=k\delta `$, where sufficiently small $`\delta >0`$ is fixed.
Applying (5.21) several times for the initial function $`\rho =\mathrm{𝟏}_I`$, we obtain, that for a Markovian moment (taking discrete values) $`t(\gamma )=k(\gamma )\delta `$
(5.22)
$$\nu _x(I)=𝔼_{h_{\gamma |_{[0,t]}}(I)}\underset{j=1}{\overset{k}{}}\frac{p(z_{j1},z_j;\delta )}{p(x_{j1},x_j;\delta )}d\nu _{\gamma (t)}(z),$$
where $`x_j=\gamma (j\delta )`$, $`z_j=h_{\gamma |_{[k\delta ,j\delta ]}}(z)`$, and the expectation is taken in the sense of the measure $`W_{x_0}`$.
Let us take, as in the previous case, for $`T=K\delta `$ sufficiently big, a Markovian moment $`\tau (\gamma )=k(\gamma )\delta `$ defined as the smallest value in the interval $`[T,2T]`$ such that for every $`t=l\delta <\tau `$
$$h_{\gamma |_{[\tau ,t]}}^{}(\gamma (\tau ))<Ce^{\beta (\tau t)}.$$
Once again, if such a moment does not exist, we take $`\tau =2T`$.
For $`T`$ sufficiently big, the probability of $`\tau <2T`$ is close to 1. Now, let us take $`I=U_{e^{\alpha T}}(x_0)`$. For most paths, as we know, the holonomy maps expand exponentially and thus $`h_{\gamma |_{[0,\tau ]}}(I)U_\epsilon (\gamma (\tau ))`$.
Note that for most of the paths starting at $`x_0`$, and for a point $`z_0`$ in the holonomy preimage $`h_{\gamma |_{[\tau ,0]}}(U_\epsilon (\gamma (\tau )))`$, the product of the quotients of the heat kernels (due to estimates analogue to these of Lemma 2.4) is bounded from below by some constant $`c_1>0`$. Let us denote
$$\begin{array}{c}N=\{\gamma \mathrm{\Gamma }_{x_0}:\tau (\gamma )<2T,h_{\gamma |_{[0,\tau ]}}(I)U_\epsilon (\gamma (\tau )),\hfill \\ \hfill z_0h_{\gamma |_{[\tau ,0]}}(U_\epsilon (\gamma (\tau )))\underset{j=1}{\overset{k}{}}\frac{p(z_{j1},z_j;\delta )}{p(x_{j1},x_j;\delta )}c_1\}\end{array}$$
Thus, the right hand side of (5.22) can be estimated as
(5.23)
$$\begin{array}{c}𝔼_{h_{\gamma |_{[0,t]}}(I)}\underset{j=1}{\overset{k}{}}\frac{p(z_{j1},z_j;\delta )}{p(x_{j1},x_j;\delta )}d\nu _{\gamma (t)}(z)\hfill \\ \hfill _N_{h_{\gamma |_{[0,t]}}(I)}\underset{j=1}{\overset{k}{}}\frac{p(z_{j1},z_j;\delta )}{p(x_{j1},x_j;\delta )}d\nu _{\gamma (t)}(z)dW_{x_0}\\ \hfill _N_{U_\epsilon (\gamma (\tau ))}c_1𝑑\nu _{\gamma (t)}(z)𝑑W_{x_0}=\\ \hfill =_Nc_1\nu _{\gamma (t)}(U_\epsilon (\gamma (\tau )))𝑑W_{x_0}W_{x_0}(N)c_1c_0.\end{array}$$
Once again we see that the measure $`\nu _{x_0}(I)`$ does not tend to 0 as $`I`$ contracts to $`x_0`$, which contradicts the fact that the measure $`\nu `$ can not have atoms.
The higher codimension case is handled in the same way by working in $`\overline{M}`$ defined in Section 5.3. The measure $`\mu `$ defines a $`\sigma `$-finite measure on $`\overline{M}`$, which harmonic in the sense of integral definition for compactly supported test function. Considering leafwise Brownian motion in $`\overline{M}`$ (with the possibility of exiting through the boundary) we see that this measure is superharmonic in the sense that $`\mu D_{}^t\mu `$, and the same is true for the conditional measures $`\nu _x`$.
Note that as in (5.23) the only trajectories used for estimates are those who arrive in the $`\epsilon `$-neighborhood of the $`\gamma (\tau )`$; hence due to exponential contraction in the inverse time they stay closer and closer to the main leaf $`_{x_0}`$. In particular, they stay in $`\overline{M}`$ for all the time in the interval $`[0,\tau ]`$.
Adding all this together, we see that the same estimates work in this case. So, the general case is handled in the same way. $`\mathrm{}`$
## 6. Acknowledgements
The authors would like to thank É. Ghys, who introduced us to the problem, and F. Ledrappier for a careful reading and for many interesting comments on this text. We also thank S. Crovisier, who explained to us the Contraction Lemma, and S. Frankel, A. Gorodetski, Yu. Ilyashenko and A. Navas for fruitful discussions. |
warning/0506/math0506272.html | ar5iv | text | # A structure theorem for quasi-Hopf comodule algebras Research partially supported by the EC programme LIEGRITS, RTN 2003, 505078, and by the bilateral project “New techniques in Hopf algebras and graded ring theory” of the Flemish and Romanian Ministries of Research. The first author was also partially supported by the programme CERES of the Romanian Ministry of Education and Research, contract no. 4-147/2004.
## Introduction
If $`H`$ is a Hopf algebra and $`B`$ is a right $`H`$-comodule algebra with the property that there exists a morphism $`v:HB`$ of right $`H`$-comodule algebras, then it is well-known that $`B`$ is isomorphic as a right $`H`$-comodule algebra to a smash product $`A\mathrm{\#}H`$, where $`A`$ is obtained as $`A=B^{co(H)}`$ and its multiplication is the restriction of the multiplication of $`B`$.
On the other hand, if $`H`$ is a quasi-bialgebra and $`A`$ is a left $`H`$-module algebra, then the smash product $`A\mathrm{\#}H`$ introduced in becomes a right $`H`$-comodule algebra and the map $`j:HA\mathrm{\#}H,j(h)=1\mathrm{\#}h`$, is a morphism of right $`H`$-comodule algebras. This raises the natural problem of checking whether, for a quasi-Hopf algebra $`H`$ and a right $`H`$-comodule algebra $`B`$ such that there exists a morphism $`v:HB`$ of right $`H`$-comodule algebras, there exists a left $`H`$-module algebra $`A`$ such that $`BA\mathrm{\#}H`$ as right $`H`$-comodule algebras. It is likely that $`A`$ appears as some sort of coinvariants of $`B`$, but it is clear that its multiplication cannot be obtained as the restriction of the one of $`B`$ (since $`B`$ is associative while in general $`A`$ is not), hence we need a different approach than in the Hopf case.
We first prove that $`B`$ becomes an object in the category $`{}_{H}{}^{}_{H}^{H}`$ of quasi-Hopf $`H`$-bimodules as introduced in . An object $`M`$ in $`{}_{H}{}^{}_{H}^{H}`$ is endowed with a projection $`E:MM`$ and a concept of coinvariants, $`M^{co(H)}=E(M)`$; applying this to $`B`$, we obtain a projection $`E:BB`$ and a subspace $`B^{co(H)}=E(B)`$. Now we define the vector space $`A=B^{co(H)}`$, with a multiplication defined by $`aa^{}=E(aa^{})`$, for all $`a,a^{}A`$. The isomorphism $`BA\mathrm{\#}H`$ follows from the structure theorem for quasi-Hopf $`H`$-bimodules, cf. , and we only have to prove that it is a morphism of right $`H`$-comodule algebras.
An application of our structure theorem is that, if we have a smash product $`A\mathrm{\#}H`$ for a quasi-Hopf algebra $`H`$, it provides a method to get $`A`$ back from $`A\mathrm{\#}H`$, as $`A(A\mathrm{\#}H)^{co(H)}`$.
## 1 Preliminaries
We work over a field $`k`$. All algebras, linear spaces etc. will be over $`k`$; unadorned $``$ means $`_k`$. Following Drinfeld , a quasi-bialgebra is a fourtuple $`(H,\mathrm{\Delta },\epsilon ,\mathrm{\Phi })`$, where $`H`$ is an associative algebra with unit, $`\mathrm{\Phi }`$ is an invertible element in $`HHH`$, and $`\mathrm{\Delta }:HHH`$ and $`\epsilon :Hk`$ are algebra homomorphisms satisfying the identities
$`(id\mathrm{\Delta })(\mathrm{\Delta }(h))=\mathrm{\Phi }(\mathrm{\Delta }id)(\mathrm{\Delta }(h))\mathrm{\Phi }^1,`$ (1.1)
$`(id\epsilon )(\mathrm{\Delta }(h))=h1,\text{ }(\epsilon id)(\mathrm{\Delta }(h))=1h,`$ (1.2)
for all $`hH`$, and $`\mathrm{\Phi }`$ has to be a normalized $`3`$-cocycle, in the sense that
$`(1\mathrm{\Phi })(id\mathrm{\Delta }id)(\mathrm{\Phi })(\mathrm{\Phi }1)=(idid\mathrm{\Delta })(\mathrm{\Phi })(\mathrm{\Delta }idid)(\mathrm{\Phi }),`$ (1.3)
$`(id\epsilon id)(\mathrm{\Phi })=111.`$ (1.4)
The identities (1.2), (1.3) and (1.4) also imply that
$$(\epsilon idid)(\mathrm{\Phi })=(idid\epsilon )(\mathrm{\Phi })=111.$$
(1.5)
The map $`\mathrm{\Delta }`$ is called the coproduct or the comultiplication, $`\epsilon `$ the counit and $`\mathrm{\Phi }`$ the reassociator. We will use the version of Sweedler’s sigma notation: $`\mathrm{\Delta }(h)=h_1h_2`$, and since $`\mathrm{\Delta }`$ is only quasi-coassociative we adopt the further convention
$`(\mathrm{\Delta }id)(\mathrm{\Delta }(h))=h_{(1,1)}h_{(1,2)}h_2,(id\mathrm{\Delta })(\mathrm{\Delta }(h))=h_1h_{(2,1)}h_{(2,2)},`$
for all $`hH`$. We will denote the tensor components of $`\mathrm{\Phi }`$ by capital letters, and those of $`\mathrm{\Phi }^1`$ by small letters, namely
$`\mathrm{\Phi }=X^1X^2X^3=T^1T^2T^3=Y^1Y^2Y^3=\mathrm{}`$
$`\mathrm{\Phi }^1=x^1x^2x^3=t^1t^2t^3=y^1y^2y^3=\mathrm{}`$
The quasi-bialgebra $`H`$ is called a quasi-Hopf algebra if there exists an anti-automorphism $`S`$ of the algebra $`H`$ and elements $`\alpha ,\beta H`$ such that, for all $`hH`$, we have:
$`S(h_1)\alpha h_2=\epsilon (h)\alpha \text{ and }h_1\beta S(h_2)=\epsilon (h)\beta ,`$ (1.6)
$`X^1\beta S(X^2)\alpha X^3=1\text{ and }S(x^1)\alpha x^2\beta S(x^3)=1.`$ (1.7)
The axioms for a quasi-Hopf algebra imply that $`\epsilon (\alpha )\epsilon (\beta )=1`$, so, by rescaling $`\alpha `$ and $`\beta `$, we may assume without loss of generality that $`\epsilon (\alpha )=\epsilon (\beta )=1`$ and $`\epsilon S=\epsilon `$.
If $`H`$ is a quasi-Hopf algebra, following , we may define the elements
$`p_R=p^1p^2=x^1x^2\beta S(x^3),q_R=q^1q^2=X^1S^1(\alpha X^3)X^2,`$ (1.8)
satisfying the relations (for all $`hH`$):
$`q_1^1p^1q_2^1p^2S(q^2)=11,q^1p_1^1S^1(p^2)q^2p_2^1=11,`$ (1.9)
$`\mathrm{\Delta }(h_1)p_R[1S(h_2)]=p_R[h1],[1S^1(h_2)]q_R\mathrm{\Delta }(h_1)=[h1]q_R.`$ (1.10)
Let us record the following easy consequence of (1.7) (for $`q=q_R=q^1q^2`$):
$`q^1\beta S(q^2)=1.`$ (1.11)
Recall from the notion of comodule algebra over a quasi-bialgebra.
###### Definition 1.1
Let $`H`$ be a quasi-bialgebra. A unital associative algebra $`B`$ is called a right $`H`$-comodule algebra if there exist an algebra morphism $`\rho :BBH`$ and an invertible element $`\mathrm{\Phi }_\rho BHH`$ such that:
$`\mathrm{\Phi }_\rho (\rho id)(\rho (b))=(id\mathrm{\Delta })(\rho (b))\mathrm{\Phi }_\rho ,\text{ }\text{ }bB\text{,}`$ (1.12)
$`(1_B\mathrm{\Phi })(id\mathrm{\Delta }id)(\mathrm{\Phi }_\rho )(\mathrm{\Phi }_\rho 1_H)=(idid\mathrm{\Delta })(\mathrm{\Phi }_\rho )(\rho idid)(\mathrm{\Phi }_\rho ),`$ (1.13)
$`(id\epsilon )\rho =id,`$ (1.14)
$`(id\epsilon id)(\mathrm{\Phi }_\rho )=(idid\epsilon )(\mathrm{\Phi }_\rho )=1_B1_H.`$ (1.15)
The first example of a right $`H`$-comodule algebra is $`H`$ itself, with $`\rho =\mathrm{\Delta }`$ and $`\mathrm{\Phi }_\rho =\mathrm{\Phi }`$. For a right $`H`$-comodule algebra $`(B,\rho ,\mathrm{\Phi }_\rho )`$ we will denote $`\rho (b)=b_{(0)}b_{(1)}`$, for all $`bB`$. If $`(B^{},\rho ^{},\mathrm{\Phi }_\rho ^{})`$ is another right $`H`$-comodule algebra, a morphism of right $`H`$-comodule algebras $`f:BB^{}`$ is an algebra map such that $`\rho ^{}f=(fid)\rho `$ and $`\mathrm{\Phi }_\rho ^{}=(fidid)(\mathrm{\Phi }_\rho )`$.
Suppose that $`(H,\mathrm{\Delta },\epsilon ,\mathrm{\Phi })`$ is a quasi-bialgebra. If $`U,V,W`$ are left $`H`$-modules, define $`a_{U,V,W}:(UV)WU(VW)`$ by
$`a_{U,V,W}((uv)w)=\mathrm{\Phi }(u(vw)).`$
The category $`{}_{H}{}^{}`$ of left $`H`$-modules becomes a monoidal category (see , for the terminology) with tensor product $``$ given via $`\mathrm{\Delta }`$, associativity constraints $`a_{U,V,W}`$, unit $`k`$ as a trivial $`H`$-module and the usual left and right unit constraints.
Let again $`H`$ be a quasi-bialgebra. We say that a $`k`$-vector space $`A`$ is a left $`H`$-module algebra if it is an algebra in the monoidal category $`{}_{H}{}^{}`$, that is $`A`$ has a multiplication and a usual unit $`1_A`$ satisfying the following conditions:
$`(aa^{^{}})a^{^{\prime \prime }}=(X^1a)[(X^2a^{^{}})(X^3a^{^{\prime \prime }})],`$ (1.16)
$`h(aa^{^{}})=(h_1a)(h_2a^{^{}}),`$ (1.17)
$`h1_A=\epsilon (h)1_A,`$ (1.18)
for all $`a,a^{^{}},a^{^{\prime \prime }}A`$ and $`hH`$, where $`haha`$ is the left $`H`$-module structure of $`A`$. Following we define the smash product $`A\mathrm{\#}H`$ as follows: as vector space $`A\mathrm{\#}H`$ is $`AH`$ (elements $`ah`$ will be written $`a\mathrm{\#}h`$) with multiplication given by
$$(a\mathrm{\#}h)(a^{^{}}\mathrm{\#}h^{^{}})=(x^1a)(x^2h_1a^{^{}})\mathrm{\#}x^3h_2h^{^{}},$$
(1.19)
for all $`a,a^{^{}}A`$, $`h,h^{^{}}H`$. Then $`A\mathrm{\#}H`$ is an associative algebra with unit $`1_A\mathrm{\#}1`$. Moreover, by , $`(A\mathrm{\#}H,\rho ,\mathrm{\Phi }_\rho )`$ becomes a right $`H`$-comodule algebra, with $`\rho :A\mathrm{\#}H(A\mathrm{\#}H)H`$, $`\rho (a\mathrm{\#}h)=(x^1a\mathrm{\#}x^2h_1)x^3h_2`$ and $`\mathrm{\Phi }_\rho =(1\mathrm{\#}X^1)X^2X^3`$. Also, it is easy to see that the map $`j:HA\mathrm{\#}H`$, $`j(h)=1\mathrm{\#}h`$, is a morphism of right $`H`$-comodule algebras.
If $`A`$, $`A^{}`$ are left $`H`$-module algebras, a map $`f:AA^{}`$ is a morphism of left $`H`$-module algebras if it is multiplicative, unital and a morphism of left $`H`$-modules.
If $`H`$ is a quasi-Hopf algebra, $`B`$ an associative algebra and $`v:HB`$ an algebra map, then, following , we can introduce on the vector space $`B`$ a left $`H`$-module algebra structure, denoted by $`B^v`$ in what follows, for which the multiplication, unit and left $`H`$-action are:
$`bb^{}=v(X^1)bv(S(x^1X^2)\alpha x^2X_1^3)b^{}v(S(x^3X_2^3)),b,b^{}B,`$ (1.20)
$`1_{B^v}=v(\beta ),`$ (1.21)
$`h_vb=v(h_1)bv(S(h_2)),hH,bB.`$ (1.22)
If $`H`$ is a quasi-Hopf algebra and $`A`$ is a left $`H`$-module algebra, define the map
$`i_0:AA\mathrm{\#}H,i_0(a)=p^1a\mathrm{\#}p^2,aA,`$ (1.23)
where $`p=p_R=p^1p^2`$ is given by (1.8). Then, by , $`i_0`$ becomes a morphism of left $`H`$-module algebras from $`A`$ to $`(A\mathrm{\#}H)^j`$.
## 2 The structure theorem
We start with a lemma which is of independent interest.
###### Lemma 2.1
Let $`H`$ be a quasi-bialgebra and $`A`$ a left $`H`$-module with a multiplication. Define a multiplication on $`AH`$ by
$`(ah)(a^{}h^{})=(x^1a)(x^2h_1a^{})x^3h_2h^{},`$ (2.1)
for all $`a,a^{}A`$ and $`h,h^{}H`$, and assume that this multiplication is associative. Then:
(i) The multiplication of $`A`$ satisfies the condition
$`(ab)c=(X^1a)((X^2b)(X^3c)),a,b,cA.`$
(ii) If moreover $`A`$ has a usual unit $`1_A`$ satisfying $`h1_A=\epsilon (h)1_A`$ for all $`hH`$, then
$`h(ab)=(h_1a)(h_2b),`$
for all $`hH`$ and $`a,bA`$, that is $`A`$ is a left $`H`$-module algebra, so the multiplication (2.1) is just the one of the smash product $`A\mathrm{\#}H`$.
Proof. (i) Let $`a,b,cA`$; then one can easily compute that in $`AH`$ we have:
$`((a1)(b1))(c1)=(y^1((x^1a)(x^2b)))(y^2x_1^3c)y^3x_2^3,`$
$`(a1)((b1)(c1))=(y^1a)(y^2((x^1b)(x^2c)))y^3x^3.`$
Since $`AH`$ is associative, these are equal; by applying $`\epsilon `$ on the second position, we obtain $`(ab)c=(X^1a)((X^2b)(X^3c))`$, q.e.d.
(ii) Let $`a,bA`$ and $`hH`$; write that $`((1_Ah)(a1))(b1)=(1_Ah)((a1)(b1))`$ in $`AH`$, then apply $`\epsilon `$ in the second position and obtain $`(h_1a)(h_2b)=h(ab)`$, q.e.d. $`\mathrm{}`$
The main ingredient for proving our structure theorem for quasi-Hopf comodule algebras will be the structure theorem for quasi-Hopf bimodules, so we recall first some facts from .
Let $`H`$ be a quasi-bialgebra and $`M`$ an $`H`$-bimodule together with an $`H`$-bimodule map $`\rho :MMH`$, with notation $`\rho (m)=m_{(0)}m_{(1)}`$ for $`mM`$ ($`\rho `$ is called a right $`H`$-coaction on $`M`$). Then $`(M,\rho )`$ is called a (right) quasi-Hopf $`H`$-bimodule if
$`(id_M\epsilon )\rho =id_M,`$ (2.2)
$`\mathrm{\Phi }(\rho id_M)(\rho (m))=(id_M\mathrm{\Delta })(\rho (m))\mathrm{\Phi },mM.`$ (2.3)
The category of right quasi-Hopf $`H`$-bimodules will be denoted by $`{}_{H}{}^{}_{H}^{H}`$ (the morphisms in the category are the $`H`$-bimodule maps intertwining the $`H`$-coactions).
If $`(V,)`$ is a left $`H`$-module, then $`VH`$ becomes a right quasi-Hopf $`H`$-bimodule with structure:
$`a(vh)b=(a_1v)a_2hb,`$ (2.4)
$`\rho _{VH}(vh)=(x^1vx^2h_1)x^3h_2,`$ (2.5)
for all $`a,b,hH`$ and $`vV`$.
Suppose now that $`H`$ is a quasi-Hopf algebra and $`(M,\rho )`$ is a right quasi-Hopf $`H`$-bimodule. Define the map $`E:MM`$, by
$`E(m)=q^1m_{(0)}\beta S(q^2m_{(1)}),mM,`$ (2.6)
where $`q=q_R=q^1q^2`$ is given by (1.8). Also, for $`hH`$ and $`mM`$, define
$`hm=E(hm).`$ (2.7)
Some properties of $`E`$ and $``$ are collected in , Proposition 3.4, for instance (for $`h,h^{}H`$ and $`mM`$): $`E^2=E`$; $`E(mh)=E(m)\epsilon (h)`$; $`hE(m)=E(hm)hm`$; $`(hh^{})m=h(h^{}m)`$; $`hE(m)=(h_1E(m))h_2`$; $`E(m_{(0)})m_{(1)}=m`$; $`E(E(m)_{(0)})E(m)_{(1)}=E(m)1`$.
Because of these properties, the following notions of coinvariants all coincide:
$`M^{co(H)}=E(M)=\{nM/E(n)=n\}=\{nM/E(n_{(0)})n_{(1)}=E(n)1\}.`$
From the above properties it follows that $`(M^{co(H)},)`$ is a left $`H`$-module.
Another description of $`M^{co(H)}`$ is (, Corollary 3.9):
$`M^{co(H)}=\{nM/\rho (n)=(x^1n)x^2x^3\}.`$
For a quasi-Hopf $`H`$-bimodule of type $`VH`$, with $`V`$$`{}_{H}{}^{}`$, we have $`(VH)^{co(H)}=V1`$ and $`E(vh)=v\epsilon (h)1`$, for all $`vV`$ and $`hH`$.
We can state now the structure theorem for quasi-Hopf $`H`$-bimodules.
###### Theorem 2.2
() Let $`H`$ be a quasi-Hopf algebra and $`M`$ a right quasi-Hopf $`H`$-bimodule. Consider $`V=M^{co(H)}`$ as a left $`H`$-module with $`H`$-action $``$ as in (2.7), and $`VH`$ as a right quasi-Hopf $`H`$-bimodule as above. Then the map
$`\nu :VHM,\nu (vh)=vh,vVandhH,`$ (2.8)
provides an isomorphism of right quasi-Hopf $`H`$-bimodules, with inverse
$`\nu ^1:MVH,\nu ^1(m)=E(m_{(0)})m_{(1)},mM.`$ (2.9)
From now on, we fix a quasi-Hopf algebra $`H`$ and a right $`H`$-comodule algebra $`(B,\rho ,\mathrm{\Phi }_\rho )`$, with notation $`\rho (b)=b_{(0)}b_{(1)}BH`$, such that there exists $`v:HB`$ a morphism of right $`H`$-comodule algebras (in particular, this implies $`\rho (v(h))=v(h_1)h_2`$, for all $`hH`$, and $`\mathrm{\Phi }_\rho =v(X^1)X^2X^3`$).
###### Lemma 2.3
$`(B,\rho )`$ becomes an object in $`{}_{H}{}^{}_{H}^{H}`$.
Proof. First, $`B`$ becomes an $`H`$-bimodule via $`v`$ (i.e. $`hbh^{}=v(h)bv(h^{})`$ for all $`h,h^{}H`$ and $`bB`$). We prove now that $`\rho :BBH`$ is an $`H`$-bimodule map. We compute:
$`\rho (hbh^{})`$ $`=`$ $`\rho (v(h)bv(h^{}))`$
$`=`$ $`\rho (v(h))\rho (b)\rho (v(h^{}))`$
$`=`$ $`(v(h_1)h_2)(b_{(0)}b_{(1)})(v(h_1^{})h_2^{})`$
$`=`$ $`v(h_1)b_{(0)}v(h_1^{})h_2b_{(1)}h_2^{}`$
$`=`$ $`h_1b_{(0)}h_1^{}h_2b_{(1)}h_2^{}`$
$`=`$ $`h\rho (b)h^{},q.e.d.`$
Obviously we have $`(id_B\epsilon )\rho =id_B`$. Finally, it is easy to see that
$`\mathrm{\Phi }(\rho id_B)(\rho (b))=(id_B\mathrm{\Delta })(\rho (b))\mathrm{\Phi },`$
because this is exactly the condition
$`\mathrm{\Phi }_\rho (\rho id_B)(\rho (b))=(id_B\mathrm{\Delta })(\rho (b))\mathrm{\Phi }_\rho `$
from the definition of a right $`H`$-comodule algebra, due to the fact that $`\mathrm{\Phi }_\rho =v(X^1)X^2X^3`$. Hence $`(B,\rho )`$ is indeed a right quasi-Hopf $`H`$-bimodule. $`\mathrm{}`$
Since $`B`$ is an object in $`{}_{H}{}^{}_{H}^{H}`$, we can consider the map $`E:BB`$, which is given by
$`E(b)=v(q^1)b_{(0)}v(\beta S(q^2b_{(1)})),bB,`$ (2.10)
where $`q=q_R=q^1q^2`$ is given by (1.8), and we can take the coinvariants
$`B^{co(H)}=E(B)=\{bB/E(b)=b\}=\{bB/E(b_{(0)})b_{(1)}=E(b)1\}.`$ (2.11)
The $`H`$-module algebra $`A`$ we are looking for will be, as a vector space, $`A=B^{co(H)}`$.
We have the $`H`$-action on $`B`$ given by $`hb=E(v(h)b)`$, which gives a left $`H`$-module structure on $`A`$. Let us note that, because of (1.11), we have $`E(1)=1`$, hence $`1A`$. By the structure theorem for quasi-Hopf $`H`$-bimodules, we know that the map
$`\mathrm{\Psi }:AHB,\mathrm{\Psi }(ah)=av(h),`$ (2.12)
is an isomorphism in $`{}_{H}{}^{}_{H}^{H}`$ (the left $`H`$-module structure of $`A`$ is $``$), with inverse
$`\mathrm{\Psi }^1:BAH,\mathrm{\Psi }^1(b)=E(b_{(0)})b_{(1)}.`$ (2.13)
Our aim will be to introduce a new multiplication on $`A`$, denoted by $``$, such that $`(A,,1,)`$ becomes a left $`H`$-module algebra and $`\mathrm{\Psi }`$ becomes an isomorphism of right $`H`$-comodule algebras between $`A\mathrm{\#}H`$ and $`B`$ (note that $`\mathrm{\Psi }`$ has the property that $`\mathrm{\Psi }j=v`$, where $`j`$ is the canonical map $`HA\mathrm{\#}H`$). We will define actually a new multiplication $``$ on the whole $`B`$, and will take its restriction to $`A`$. Namely, for all $`b,b^{}B`$, define
$`bb^{}=E(bb^{}).`$ (2.14)
Since $`A=E(B)`$, $``$ restricts to a multiplication on $`A`$. Since for $`aA`$ we have $`E(a)=a`$, we obtain $`a1=1a=E(a)=a`$, hence $`1`$ is a unit for $`(A,)`$. Let now $`hH`$; we compute:
$`h1`$ $`=`$ $`E(v(h))`$
$`=`$ $`v(q^1)v(h)_{(0)}v(\beta S(q^2v(h)_{(1)}))`$
$`=`$ $`v(q^1)v(h_1)v(\beta S(q^2h_2))`$
$`=`$ $`v(q^1h_1\beta S(h_2)S(q^2))`$
$`(\text{1.6})`$ $`=`$ $`v(q^1\beta S(q^2))\epsilon (h)`$
$`(\text{1.11})`$ $`=`$ $`\epsilon (h)1.`$
In view of Lemma 2.1, in order to get that $`(A,,1,)`$ is a left $`H`$-module algebra, it is enough to prove that the multiplication defined on $`AH`$ by
$`(ah)(a^{}h^{})=(x^1a)(x^2h_1a^{})x^3h_2h^{}`$
is associative. Since $`\mathrm{\Psi }:AHB`$ is bijective and $`B`$ is associative, it is enough to prove that $`\mathrm{\Psi }`$ is multiplicative, that is, for all $`a,a^{}A`$ and $`h,h^{}H`$:
$`\mathrm{\Psi }((x^1a)(x^2h_1a^{})x^3h_2h^{})=\mathrm{\Psi }(ah)\mathrm{\Psi }(a^{}h^{}).`$
We prove first a relation that will be used in the proof of the multiplicativity of $`\mathrm{\Psi }`$.
###### Lemma 2.4
Let $`H`$ be a quasi-Hopf algebra; then we have:
$`q_1^1t^1x^1q_{(2,1)}^1t_1^2z^1x^2q_{(2,2)}^1t_2^2z^2\beta S(q^2t^3z^3)x^3=111,`$ (2.15)
where $`q=q_R=q^1q^2`$ is given by (1.8).
Proof. We will use also the element $`p=p_R=p^1p^2`$ given by (1.8). We compute:
$`q_1^1t^1x^1q_{(2,1)}^1t_1^2z^1x^2q_{(2,2)}^1t_2^2z^2\beta S(q^2t^3z^3)x^3`$
$`(\text{1.3})`$ $`=`$ $`q_1^1Z^1t_1^1y^1x^1q_{(2,1)}^1Z^2t_2^1y^2x^2q_{(2,2)}^1Z^3t^2y_1^3\beta S(q^2t^3y_2^3)x^3`$
$`(\text{1.6}),(\text{1.5})`$ $`=`$ $`q_1^1Z^1t_1^1x^1q_{(2,1)}^1Z^2t_2^1x^2q_{(2,2)}^1Z^3t^2\beta S(q^2t^3)x^3`$
$`(\text{1.1}),(\text{1.8})`$ $`=`$ $`Z^1q_{(1,1)}^1p_1^1x^1Z^2q_{(1,2)}^1p_2^1x^2Z^3q_2^1p^2S(q^2)x^3`$
$`(\text{1.9})`$ $`=`$ $`Z^1x^1Z^2x^2Z^3x^3`$
$`=`$ $`111,`$
and the relation is proved. $`\mathrm{}`$
We will need two of the general properties of the map $`E`$ on a right quasi-Hopf $`H`$-bimodule $`M`$ recalled before, which for $`M=B`$ become:
$`v(h)a=(h_1a)v(h_2),hH,aA,`$ (2.16)
and, if $`aB`$, then
$`aAa_{(0)}a_{(1)}=(x^1a)v(x^2)x^3.`$ (2.17)
We need also a more explicit formula for $`aa^{}`$, if $`a,a^{}A`$. We can write:
$`aa^{}`$ $`=`$ $`E(aa^{})`$
$`=`$ $`v(q^1)a_{(0)}a_{(0)}^{}v(\beta S(q^2a_{(1)}a_{(1)}^{}))`$
$`(\text{2.17})`$ $`=`$ $`v(q^1)(t^1a)v(t^2)(z^1a^{})v(z^2\beta S(q^2t^3z^3)).`$
We can finally prove that $`\mathrm{\Psi }`$ is multiplicative. We compute (for $`a,a^{}A`$ and $`h,h^{}H`$):
$`\mathrm{\Psi }(ah)\mathrm{\Psi }(a^{}h^{})`$ $`=`$ $`av(h)a^{}v(h^{})`$
$`(\text{2.16})`$ $`=`$ $`a(h_1a^{})v(h_2h^{}),`$
$`\mathrm{\Psi }((x^1a)(x^2h_1a^{})x^3h_2h^{})`$
$`=`$ $`(x^1a)(x^2h_1a^{})v(x^3h_2h^{})`$
$`=`$ $`v(q^1)(t^1x^1a)v(t^2)(z^1x^2h_1a^{})v(z^2\beta S(q^2t^3z^3)x^3h_2h^{})`$
$`(\text{2.16})`$ $`=`$ $`(q_1^1t^1x^1a)v(q_2^1t^2)(z^1x^2h_1a^{})v(z^2\beta S(q^2t^3z^3)x^3h_2h^{})`$
$`(\text{2.16})`$ $`=`$ $`(q_1^1t^1x^1a)(q_{(2,1)}^1t_1^2z^1x^2h_1a^{})v(q_{(2,2)}^1t_2^2z^2\beta S(q^2t^3z^3)x^3h_2h^{})`$
$`(\text{2.15})`$ $`=`$ $`a(h_1a^{})v(h_2h^{}),q.e.d.`$
Since obviously we have $`\mathrm{\Psi }(11)=1`$, now we have that $`(A,,1,)`$ is a left $`H`$-module algebra and $`\mathrm{\Psi }:A\mathrm{\#}HB`$ is an algebra isomorphism. Using (2.17), the fact that $`\rho (v(h))=v(h_1)h_2`$ and the formula $`\rho _{A\mathrm{\#}H}(a\mathrm{\#}h)=(x^1a\mathrm{\#}x^2h_1)x^3h_2`$, one can easily see that $`\rho _B\mathrm{\Psi }=(\mathrm{\Psi }id)\rho _{A\mathrm{\#}H}`$. Moreover, since $`\mathrm{\Phi }_{A\mathrm{\#}H}=(1\mathrm{\#}X^1)X^2X^3`$ and $`\mathrm{\Psi }(1\mathrm{\#}X^1)X^2X^3=v(X^1)X^2X^3=\mathrm{\Phi }_B`$, we conclude that $`\mathrm{\Psi }`$ is an isomorphism of right $`H`$-comodule algebras. Hence, we have proved the desired structure theorem:
###### Theorem 2.5
Let $`H`$ be a quasi-Hopf algebra and $`B`$ a right $`H`$-comodule algebra such that there exists $`v:HB`$ a morphism of right $`H`$-comodule algebras. Then there exists a left $`H`$-module algebra $`A`$ (whose structure is described above) such that $`BA\mathrm{\#}H`$ as right $`H`$-comodule algebras.
Let now $`H`$, $`B`$ and $`v:HB`$ be as above. Since $`B`$ is an associative algebra and $`v`$ is an algebra map, we can consider the left $`H`$-module algebra $`B^v`$ as in the Preliminaries.
###### Proposition 2.6
With notation as above, the map
$`\theta :AB^v,\theta (a)=(p^1a)v(p^2),`$
where $`p=p_R=p^1p^2`$ is given by (1.8), is an injective morphism of left $`H`$-module algebras.
Proof. Since $`\mathrm{\Psi }:A\mathrm{\#}HB`$ is an algebra map satisfying $`\mathrm{\Psi }j=v`$, by , Lemma 4.1. it follows that $`\mathrm{\Psi }:(A\mathrm{\#}H)^jB^v`$ is a morphism of left $`H`$-module algebras. We know from the Preliminaries that the map $`i_0:A(A\mathrm{\#}H)^j`$, $`i_0(a)=p^1ap^2`$, is also a morphism of left $`H`$-module algebras, and one can see that actually $`\theta =\mathrm{\Psi }i_0`$, hence $`\theta `$ is indeed a morphism of left $`H`$-module algebras, and it is injective since $`i_0`$ is injective and $`\mathrm{\Psi }`$ is bijective.
Note that in the Hopf case $`\theta `$ is simply the inclusion of $`A`$ into $`B^v`$. $`\mathrm{}`$
###### Remark 2.7
Let $`H`$ be a quasi-Hopf algebra, $`A`$ a left $`H`$-module algebra and $`B=A\mathrm{\#}H`$; then, for this $`B`$ together with the canonical map $`j:HA\mathrm{\#}H`$, one can show that the map $`E`$ is given by $`E(a\mathrm{\#}h)=\epsilon (h)(a\mathrm{\#}1)`$, and so $`B^{co(H)}=A\mathrm{\#}1`$, with multiplication and $`H`$-action:
$`(a\mathrm{\#}1)(a^{}\mathrm{\#}1)=aa^{}\mathrm{\#}1,a,a^{}A,`$
$`h(a\mathrm{\#}1)=ha\mathrm{\#}1,hHandaA,`$
that is $`B^{co(H)}A`$ as left $`H`$-module algebras. Hence, the structure theorem allows to recover the structure of $`A`$ from the one of $`A\mathrm{\#}H`$. |
warning/0506/cond-mat0506470.html | ar5iv | text | # Optical spectroscopy of trivalent chromium in sol-gel lithium niobate
## Abstract
We report on the characterization of sol-gel derived lithium niobate via trivalent chromium probe ions, a study that is motivated by recent reports on the synthesis of high quality sol-gel lithium niobate (LiNbO<sub>3</sub>). In order to assess the quality of sol-gel derived LiNbO<sub>3</sub>, we incorporate Cr<sup>3+</sup> during the hydrolysis stage of the sol-gel process. A comparison of the Cr<sup>3+</sup> emission and photo-excitation data on both sol-gel and melt-grown LiNbO<sub>3</sub> shows that the sol-gel derived material is highly stoichiometric.
The need for sol-gel derived ferroelectric coatings and films for electro-optical and piezo-electric applications is driving the research for high quality sol-gel lithium niobate (LiNbO<sub>3</sub>)Shin-ichi Hirano et al. (2002); Cheng et al. (2001). Recent reports concentrate on optimizing the sol-gel process with the aim of producing high quality films with preferred orientation. X-ray diffraction and ICP have been the primary techniques to assess stoichiometry of the samples. To our knowledge, optical probes, such as Cr<sup>3+</sup> ions, have not been used to characterize the sol-gel samples, although this approach has been very successful to show fundamental differences between congruent and stoichiometric melt-grown LiNbO<sub>3</sub> Salley et al. (2000). The Cr<sup>3+</sup> ion is particulary suited to study the quality of LiNbO<sub>3</sub>, because it is arguably the most studied optical impurity Henderson and Imbusch (1989), has a spectrum that is very sensitive to the possible lattice sites and disorder in LiNbO<sub>3</sub>, and can be readily incorporated in an early stage of the sol-gel process. In addition, chromium doped lithium niobate has stimulated much interest due to the broadband near infrared luminescence of the chromium ions, with the ultimate goal of developing a tuneable laser in the visible range through frequency doubling within the active medium.
In the current work, we study the optical Cr<sup>3+</sup> transitions between the $`{}_{}{}^{4}A_{2}^{}`$ groundstate and the lowest $`{}_{}{}^{4}T_{2}^{}`$ and $`{}_{}{}^{2}E`$ excited electronic states. These transitions are characterized by broad absorption and emission bands for the spin allowed transition, and spectrally narrow lines for the spin-forbidden transition, respectively. These transitions provide information on the stoichiometry and disorder of the sol-gel material, and allow a comparison to the detailed study of LiNbO<sub>3</sub> samples grown from melt.
For the sample synthesis, we combine 1 M solutions of lithium and niobium ethoxides in ethanol such that the ratio of Li:Nb is 1:1. An appropriate amount of Cr(NO<sub>3</sub>)<sub>3</sub> is then dissolved in water (in this study the Cr:Li ratio was 0.001:1). Under constant stirring, the water-chromium solution is slowly added
into the mixed ethoxide. After an initial release of ethanol, the solution is capped and aged for several days. The aged solution is then dried at room temperature until a fine white powder has formed. The powder samples are finally heated in air to 800 <sup>o</sup>C for two hours to remove excess water and organic complexes. The calcined powders are characterized by x-ray diffraction utilizing Cu K-alpha radiation; a typical pattern is shown in Fig. 1. All observed diffraction peaks can be attributed to lithium niobate, and we find no evidence for alternate phases such as Li<sub>0.88</sub>H<sub>0.12</sub>NbO<sub>3</sub>.
The optical experiments are carried out with the powder sample mounted in an Oxford Instruments temperature variable cryostat. For emission measurements the sample is excited with a HeNe laser at 632.8 nm. The broadband emission was collected with reflective optics and analyzed with a Bruker 66v FTIR spectrometer. Photoexcitation measurements were performed by mounting the cryostat directly into the sample space of a modified Cary 14 spectrophotometer. In both the broadband emission and photoexcitation experiments, the emitted light was collected through appropriate filters and imaged onto a liquid-nitrogen cooled InGaAs detector. For the photoexcitation, a lock-in amplifier was used for signal-to-noise enhancement.
Experimental results for the sol-gel LiNbO<sub>3</sub>:0.1%Cr<sup>3+</sup> sample are shown in Fig. 2. The excitation spectrum (Fig.2, left) observed at $`\lambda _{det}`$ 900 nm shows the characteristic <sup>4</sup>A<sub>2</sub> $``$ <sup>4</sup>T<sub>2</sub> and <sup>4</sup>A<sub>2</sub> $``$ <sup>4</sup>T<sub>1</sub> broad absorption bands around 655 nm and 480 nm, respectively, as well as the spin-forbidden <sup>4</sup>A<sub>2</sub> $``$ <sup>2</sup>E transition around 726 nm. The emission spectrum (Fig. 2, right) shows the <sup>4</sup>T<sub>2</sub> $``$ <sup>4</sup>A<sub>2</sub> emission band under HeNe laser excitation ($`\lambda _{ex}`$ = 632.8 nm).
To interpret these results, they are compared to spectra obtained on standard melt-grown LiNbO<sub>3</sub>:Cr<sup>3+</sup>. Lithium niobate crystals pulled from a melt with equal amounts of lithium and niobium crystallize with a Li 3% deficit Byer et al. (1970), and are referred to as congruent samples. When chromium ions are added to these melts, the trivalent impurities preferentially occupy the crystallographic site of lithium vacancies, in the following referred to as Cr\[Li\]. The crystal field leads to strong absorption into the <sup>4</sup>T<sub>2</sub> and <sup>4</sup>T<sub>1</sub> levels peaking at 654 nm and 478 nm, respectively, giving the crystals a green color. A shift in the optical spectra is observed when LiNbO<sub>3</sub>:Cr<sup>3+</sup> is codoped with more that 4.5% Mg<sup>2+</sup>, resulting in Cr<sup>3+</sup> impurities preferentially occupying niobium sites Macfarlane et al. (1995). The resulting change in the crystal field shifts the <sup>4</sup>T<sub>2</sub> and <sup>4</sup>T<sub>1</sub> absorption peaks to 715 nm and 540 nm, resulting in a crystal with a red color. New crystal growth techniques, namely high-temperature top-seeded solution growth methods (HTTSSG), enabled the production of LiNbO<sub>3</sub> crystals that are stochiometric (i.e. no Li or Nb deficit)Polgar et al. (1997). These samples produce spectra with narrower linewidths enabling more precise optical spectroscopy and better characterization of the various chromium sites Salley et al. (2000), moreover, the chromium ions in these stoichiometric melt-grown single crystals are found to occupy both Li and Nb sites Salley (2000).
In contrast to the Cr\[Li\] spectrum that dominates in congruent LiNbO<sub>3</sub> crystals, the sol-gel produced material has a significant contribution to the luminescence at excitation wavelengths beyond 750 nm, as evident in a shoulder in the excitation spectrum extending from about 750nm to 800 nm in Fig.2. Fig.3 shows an expanded excitation spectrum in this region (dots). This shoulder can be attributed to chromium ions on niobium sites, Cr\[Nb\]. To quantify this contribution to the excitation spectrum, we measured the spectra of LiNbO<sub>3</sub>:Cr<sup>3+</sup> crystals with and without magnesium codoping (6%).
The inset of Fig.3 shows the absorption for the single crystal samples of LiNbO<sub>3</sub>:Cr (solid) and LiNbO<sub>3</sub>:Cr:Mg (dashed) in the 700 nm to 850 nm range, clearly demonstrating the difference of the <sup>4</sup>T<sub>2</sub> onset for the two sites. The fit to the sol-gel data (Fig. 3, main plot, line) was found by varying the contributions to the impurity absorption for each site while keeping the total constant, i.e $`\alpha _{solgel}`$ = C$`{}_{Li}{}^{}\alpha _{Cr[Li]}^{}`$\+ C$`{}_{Nb}{}^{}\alpha _{Cr[Nb]}^{}`$, with C<sub>Li</sub> \+ C<sub>Nb</sub> = 1. The fit shown in Fig.3 (solid line) represents a C<sub>Li</sub>=0.66 contribution from the Cr\[Li\] and a C<sub>Nb</sub> = 0.34 contribution from Cr\[Nb\]. This approach is justified because our single crystal reference samples contained the same amount of chromium ions (0.25%). A careful inspection of LiNbO<sub>3</sub>:Cr:Mg (insert Fig.3, dashed line) shows a small peak at 727 nm (<sup>4</sup>A$`{}_{2}{}^{}_{}^{2}`$E transition), due to Cr\[Li\]. Taking this correction into account, we arrive at relative Cr\[Li\] and Cr\[Nb\] concentrations of 0.64 and 0.36, respectively.
The substantial occupation of the niobium site by chromium ions is in marked difference to melt-grown congruent samples, where the chromium ions occupy lithium sites exclusively. In highly stoichiometric LiNbO<sub>3</sub> samples, grown by the HTTSSG method Polgar et al. (1997), we also find chromium ions on both the lithium and niobium sites Salley (2000), thus our results confirm reports that the sol-gel process yields highly stoichiometric samples without the intrinsic Li deficit, which is characteristic for congruent melt-grown LiNbO<sub>3</sub>. In stoichiometric single crystal samples of highest quality, the zero-phonon line of the <sup>4</sup>A$`{}_{2}{}^{}_{}^{4}`$T<sub>2</sub> can be observed at low temperature. In our sample, this transition is broadened, which is indicative of lattice defects. Our sample also shows the optical signature of so-called high field sites, which is due to Cr<sup>3+</sup> ions in a strong crystal field, resulting in the <sup>2</sup>E level lying below the <sup>4</sup>T<sub>2</sub> level, characterized by sharp emission lines at low temperatures, and which are associated with defects in the crystal. Thus, while our results indicate that the sol-gel samples do not have the intrinsic lithium deficit of congruent melt-grown material, the overall quality of the material does not yet approach that of the highest quality single crystal material. However, with the optical absorption and emission spectra of trivalent chromium being a clear indicator of the crystal quality, systematic studies are possible to improve the quality of sol-gel derived LiNbO<sub>3</sub>.
In summary, we have shown Cr<sup>3+</sup> ions are sensitive optical probes of the quality of sol-gel derived LiNbO<sub>3</sub>. Our experiments confirm that sol-gel produced lithium niobate is highly stoichiometric and does not contain the intrinsic Li deficit of congruent LiNbO<sub>3</sub>. Our results also show that our sol-gel material contains more defects than high quality stoichiometric LiNbO<sub>3</sub>. With the aid of the chromium optical probe ions, experiments are under way to improve the quality of sol-gel LiNbO<sub>3</sub> and related ferroelectric materials. |
warning/0506/astro-ph0506100.html | ar5iv | text | # Supernova Neutrino-Effects on R-Process Nucleosynthesis in Black Hole Formation
## 1 Introduction
In the past eighteen years after SN1987A emerged in the Large Magellanic Cloud, a number of interesting scenarios have been proposed to explain why a pulsar has not been discovered at the supernova (hereafter, SN) remnant. One of the most viable scenarios is the possibility that the pulsation is too weak for emitted X-rays to penetrate through the ambient gas clouds surrounding the SN remnant. Another interesting possibility is that the proto-neutron star (hereafter, NS) was first formed after the core-collapse, emitting thermal neutrinos in twelve seconds during the Kelvin-Helmholtz cooling phase, then became a black hole (hereafter, BH) after deleptonization (Brown & Bethe, 1994).
BH formation in core-collapse SNe leads to an interesting theoretical study. If BH formation happens when the flux of neutrinos is still high, abrupt termination within $`0.5`$ ms would emerge as a sharp cutoff of neutrino luminosity at some time after the core bounce. This feature allows celestial model-independent time-of-flight mass tests for the three light neutrino families (Beacom et al., 2001). Beacom et al. (2001) discussed detectable neutrino mass sensitivity that depends strongly on the timing, i.e. how early BH formation occurs after core collapse.
Motivated by these theoretical studies, we consider the possibility that the r-process nucleosynthesis could provide an independent observable signature showing BH formation during high neutrino luminosity epoch. All flavors of neutrinos and antineutrinos from the neutrino-spheres inside the proto-NS are thought to play at least two essential roles in successful SN explosions in “delayed explosion” model. First, successive interactions between intensive flux of neutrinos and materials collapsing into the proto-NS deposit neutrino energy into the ejecta and thus revive the shock propagation leading to a successful breakout through the iron core (Wilson, 1985; Bethe & Wilson, 1985). Second, in such a neutrino-powered SN explosion mechanism, the atmosphere of the proto-NS is heated by neutrinos at high entropy $`s/k`$ = $`100400`$ to form a “hot bubble” flowing out rapidly behind the shock, which is called the neutrino-driven wind. This is a viable candidate site for the r-process (Woosley et al., 1994).
The evolution of the neutrino-driven wind begins from high temperatures about $`10^{10}`$K at high entropy $`s/k`$ = $`100400`$, thus the system is in nuclear statistical equilibrium (NSE) and favors free neutrons, protons, and some amount of $`\alpha `$-particles. As the temperature drops below about $`kT0.5`$ MeV, fast charged particle reactions, which are responsible for interconverting protons into $`\alpha `$-particles and converting $`\alpha `$-particles into composite nuclei, very quickly accumulate “seeds” that have large masses $`70A120`$ ($`\alpha `$-process). This lasts until the temperature drops to $`kT0.5/e0.2`$ MeV when the charged particle reactions freeze out ($`\alpha `$-rich freezeout). Below this temperature, only neutron-capture flow goes on, followed by $`\beta `$-decays, and the r-process occurs until the neutrons are exhausted (r-process and freezeout). In Figure 1 we display the evolution of temperature (top panel) and chemical composition of the light elements (bottom panel), indicating the sequence of nuclear reaction processes, $`i.e.`$ the NSE at $`tt_\alpha `$, the $`\alpha `$-process at $`t_\alpha tt_\mathrm{n}`$, and the r-process at $`t_\mathrm{n}tt_\mathrm{f}`$. Here, we define the beginning of the $`\alpha `$-process at $`t=t_\alpha `$ when $`kT=0.5`$ MeV, the $`\alpha `$-rich freezeout at $`t=t_\mathrm{n}`$ when $`kT=0.5/e0.2`$ MeV, and the r-process freezeout at $`t=t_\mathrm{f}`$. We also define the dynamical expansion time of the wind, $`\tau _{dyn}`$, as the e-fold decay time of the temperature from $`kT=0.5`$ MeV (Qian & Woosley, 1996), i.e. $`\tau _{dyn}=t_\alpha t_\mathrm{n}`$. A short dynamical time can suppress the overproduction of seed nuclei, leaving plenty of free neutrons for the subsequent neutron-capture flow.
The important neutrino reactions during the nucleosynthesis are
$`\nu _e+{}_{N}{}^{Z}A{}_{N1}{}^{Z+1}A+e^{},`$ (1)
$`\overline{\nu _e}+{}_{N}{}^{Z}A{}_{N+1}{}^{Z1}A+e^+,`$ (2)
$`\nu _\mathrm{x}(\overline{\nu _\mathrm{x}})+_N^ZA\left[\begin{array}{c}^{Z1}_NA+\mathrm{p}\\ ^Z_{N1}A+\mathrm{n}\end{array}\right]+\nu _\mathrm{x}^{^{}}(\overline{\nu _\mathrm{x}^{^{}}}),`$ (5)
where x= $`\mu `$, and $`\tau `$ are the neutrino flavors, and $`{}_{N}{}^{Z}A`$ is the nucleus with proton number Z and neutron number N. In particular the charged-current reactions that determine the initial neutron-to-proton ratio are
$`\nu _e+\mathrm{n}\mathrm{p}+\mathrm{e}^{},`$ (6)
$`\overline{\nu _e}+\mathrm{p}\mathrm{n}+\mathrm{e}^+.`$ (7)
The neutron-to-proton ratio in the weak equilibrium satisfies (Qian & Woosley, 1996),
$`Y_e={\displaystyle \frac{\mathrm{p}}{\mathrm{n}+\mathrm{p}}}\left(1+{\displaystyle \frac{L_{\overline{\nu _e}}}{L_{\nu _e}}}\times {\displaystyle \frac{ϵ_{\overline{\nu _e}}2\mathrm{\Delta }+1.2\mathrm{\Delta }^2/ϵ_{\overline{\nu _e}}}{ϵ_{\nu _e}+2\mathrm{\Delta }+1.2\mathrm{\Delta }^2/ϵ_{\nu _e}}}\right),`$ (8)
where $`Y_e`$ is the electron fraction, $`L_\nu `$ is the neutrino luminosity of each species, $`ϵ_\nu `$ is the average energy proportional to $`T_\nu `$, and $`\mathrm{\Delta }`$ is the mass difference between proton and neutron. In this paper we take $`kT_{\nu _e}=11.0`$ MeV, $`kT_{\overline{\nu _e}}=19.0`$ MeV, and $`kT_{\nu _\mathrm{x},\overline{\nu _\mathrm{x}}}=25.0`$ MeV from Woosley et al. (1994) and Qian et al. (1997). Because of the hierarchy of different neutrino flavors $`T_\nu <T_{\overline{\nu _e}}<T_{\nu _\mathrm{x}(\overline{\nu _\mathrm{x}})}`$, which is the consequence of the different diffusion length scales and the decreasing temperature with increasing radius, this $`Y_e`$ value is less than 0.5 when $`L_{\nu _e}=L_{\overline{\nu _e}}`$, as indicated by the numerical simulations of the neutrino transfer (Woosley et al., 1994). Therefore, the whole nucleosynthesis sequence described above and in Figure 1 occurs in the neutron-rich environment with $`Y_e=\frac{\mathrm{p}}{\mathrm{n}+\mathrm{p}}<0.5`$.
Since the timescale of the $`\alpha `$ production is shorter than those of the reactions (4) and (5), protons are locked up into $`\alpha `$ particles until they are exhausted. Therefore the reaction (5) does not occur at later times. Thus, the nuclei are neutron rich and the neutrons in the nuclei are degenerate up to higher energy levels than those of the protons. Hence, the reaction (1), which is fundamentally equivalent to the reaction (4), is energetically favorable. In contrast, the reaction (2) can not occur because of the Pauli exclusion principle. Both reactions (1) and (4) have an important role in decreasing the neutron abundance. If the BH is formed and the neutrino luminosity is cut off, neutron abundance is kept high. This may lead to an efficient production of r-process elements because the neutron-to-seed ratio are large under such high neutron abundance conditions (Terasawa et al., 2001). We note that among all neutral-current reactions (3), the spallation of $`\alpha `$-particles;
$`\nu _\mathrm{x}(\overline{\nu _\mathrm{x}})+{}_{}{}^{4}\mathrm{He}\left[\begin{array}{c}^3\mathrm{H}+\mathrm{p}\\ ^3\mathrm{He}+\mathrm{n}\end{array}\right]+\nu _\mathrm{x}^{^{}}(\overline{\nu _\mathrm{x}^{^{}}}),`$ (11)
is most important to accelerate the termination of charged particle reactions during the $`\alpha `$-process by accumulating more seed nuclei and leaving less free neutrons at the end of the $`\alpha `$-process at $`tt_\mathrm{n}`$ (Meyer, 1995). This also strongly affects the r-process nucleosynthesis.
As such, the neutrino cutoff may affect nucleosynthesis because of robust emission of neutrinos before the cutoff time $`t_{cut}`$. Produced r-process abundance pattern may differ from that produced in the SNe leading to the formation of NS. This difference might be important to understand the mechanism of BH formation, which remains an open question.
In our calculations we ignore the effects of neutrino-flavor mixing (Qian and Fuller, 1995). Neutrino mixing can significantly alter the r-process nucleosynthesis yields (for a recent discussion see Balantekin and Yuksel (2005)).
This paper will be organized in the following way. In sect. 2, we first review various scenarios of BH formation, and then discuss the conditions under which the neutrino luminosity is cut off in the manner which we will discuss in the present article. In sect. 3, we describe flow models of neutrino-driven wind in the core-collapse SNe which we adopt in our numerical studies. Several assumptions and approximations are explained. We also explain the nuclear reaction network code in this section. In sect. 5, we show the calculated results on the r-process nucleosynthesis and discuss how different r-element abundances we would obtain, depending on the BH vs. NS formation after the core-collapse. Finally, in sect. 6, we summarize the present paper and discuss the implications of the spectroscopic observations of actinide elements, <sup>232</sup>Th and <sup>235,238</sup>U, in metal-deficient halo stars for testing our results.
## 2 Black Hole Formation and Neutrino Cutoff
We calculate the r-process nucleosynthesis in the SNe where BH is formed. Since we do not know when the BH is formed, we treat the time $`t_{cut}`$, at which the neutrino is cut off, as a parameter. We adopt the constant luminosity, $`L_\nu =10^{51}\mathrm{ergs}^1`$ at $`tt_{cut}`$, and assume that neutrino luminosity is zero after the neutrino cutoff, $`L_\nu =0`$ at $`t_{cut}<t`$. We investigate the dependence of the r-process nucleosynthesis on the value of $`t_{cut}`$.
Stars more massive than 8$`M_{}`$ are known to culminate their main sequence lives as core-collapse SNe. In a standard SN explosion where the progenitor star has a mass $``$ 20$`M_{}`$, in general, a NS is formed as a remnant in the center. In heavier progenitor stars, a BH can be formed instead of NS. From a cosmological point of view, the recent discovery of an extremely metal-deficient star suggests that such a supermassive star that induces a BH existed at least in the early universe (Frebel et al., 2005; Iwamoto et al., 2005).
There are so far two different scenarios for BH formation in SN explosions. One of them is a collapsar model. This model is based on a completely different picture from the standard SN. In this model a BH is thought to be formed immediately after a core collapse of SNe. It is phenomenologically possible when the mass of the progenitor star is heavy enough ($`e.g.`$ 40$`M_{}`$) to produce a BH. An accretion disk is formed around a BH by pulling materials into the central region and strong X-rays or $`\gamma `$-rays are emitted. This scenario is sometimes invoked to explain $`\gamma `$-ray bursts (MacFadyen & Woosley, 1999).
A second scenario is based on a standard core collapse SN explosion. If the proto-NS mass exceeds the maximum NS mass, then a BH can be formed during the SN explosion. Here the maximum mass is thought to be 2.2$`M_{}`$ (Beacom et al., 2001; Akmal et al., 1998). This scenario can be useful for the progenitor mass region heavier than that of a standard SN and lighter than that of a collapsar ($`e.g.`$ between 20$`M_{}`$ 40$`M_{}`$). It is the scenario for which the neutrino cutoff effect by the BH formation is obvious, since there can exist a proto-NS $`i.e.`$ neutrino flux from the central neutrino sphere for a short time until a BH is formed. This situation is different from the collapsar model.
## 3 Flow Models of Neutrino-Driven Wind
### 3.1 Background
The flow dynamics in non-spherical SN explosion is complicated in the case of collapsars (MacFadyen & Woosley, 1999) or hypernovae (Maeda & Nomoto, 2003) which are associated with BH formation. The flow in which BH is formed might be different from the standard SN flow leading to a NS formation. We however assume similar models for both flows for the following reasons. The first reason is that in this work we are mainly interested in the consequence of the neutrino-cutoff, i.e. how it destroys neutron richness and affects the final abundance of r-process elements. Since the r-process condition in Type-II SN explosion leaving NS as a remnant has been studied very well (Woosley et al., 1994; Witti et al., 1994; Otsuki et al., 2000), it is effective to extensively study the r-process in BH formation in similar models by tuning flow parameters of the neutrino-driven wind.
The second reason is that the luminosity stays at high value of $`10^{51}10^{52}`$ erg$`\mathrm{s}^1`$ for each neutrino species before the neutrino-cutoff at 1-2 s in either models of BH formation as discussed in the previous section. As we will discuss below, neutrino heating energy is so efficiently deposited in very short period $`3`$ ms (Otsuki et al., 2000) that the hot bubble may form even in the BH formation.
The third reason is that we have not yet obtained a realistic theoretical simulation of a successful SN explosion except for the model of Wilson (1985). It may be more difficult to simulate BH formation which needs at least two dimensional numerical analysis including the effects of general relativity in order to describe both the dynamics of the accretion disk and jet formation as well as the inner core collapse into BH.
For these reasons we use the approximation of the neutrino-driven wind which is suitable for our purpose to study the neutrino-cutoff effects on the r-process nucleosynthesis.
### 3.2 Steady-State Flow Model
We adopt the spherical steady-state flow model for the neutrino-driven wind (Qian & Woosley, 1996; Takahashi & Janka, 1997; Otsuki et al., 2000; Wanajo et al., 2001). This flow model is one which leads to a successful r-process. Even though the entropy per baryon is moderately low, $`s/k`$ 100-300, the r-process can occur in this neutrino-driven wind when the dynamical expansion timescale becomes much shorter than the collision timescale of neutrino-nucleus interactions. For the present application, such hydrodynamic flow can be approximated (Otsuki et al., 2003) by solving the following non-relativistic equations:
$`4\pi r^2\rho v=\dot{M},`$ (12)
$`{\displaystyle \frac{1}{2}}v^2{\displaystyle \frac{GM}{r}}+N_As_{rad}\mathrm{kT}=E,`$ (13)
$`s_{rad}={\displaystyle \frac{11\pi ^2}{45\rho N_A}}\left({\displaystyle \frac{\mathrm{k}T}{\mathrm{}c}}\right)^3,`$ (14)
where $`\dot{M}`$ is the rate at which matter is ejected by neutrino heating on the surface of the proto-NS. In Eq. (9), the total energy $`E`$ is fixed by the boundary condition on the asymptotic temperature, $`\mathrm{T}_a`$;
$`E=N_As_{rad}\mathrm{k}\mathrm{T}_a,`$ (15)
and we take into account only photons and $`e^\pm `$ pairs in the estimate of the entropy per baryon, $`s_{\mathrm{rad}}`$, in Eq. (10). For simplicity, in the present work, we utilize an adiabatic, constant-entropy wind rather than computing neutrino heating explicitly (Otsuki et al., 2000). However, we include both charged- and neutral-current interactions between neutrinos and nuclei in the nucleosynthesis. A constant neutrino luminosity $`L_\nu =10^{51}\mathrm{ergs}^1`$ for each neutrino species is also adopted. This model has four parameters, which are the NS mass, entropy, boundary temperature, and mass loss rate.
For the purpose of investigating many different flow models, we use the exponential model (Otsuki et al., 2003) which satisfies
$`\tau _{dyn}^1={\displaystyle \frac{1}{TT_a}}{\displaystyle \frac{\mathrm{d}T}{\mathrm{d}t}},`$ (16)
where $`T_a`$ is the asymptotic temperature as defined in Eq. (11) so that the solution can be a good approximation to the exact solution of the set of Eqs. (8) - (11). Under the approximate condition that the entropy ($`\mathrm{s}=(11\pi ^2/45\mathrm{N}_\mathrm{A})(T^3/\rho )`$) is constant, temperature and density are given for a fixed $`\tau _{dyn}`$. Among many different flow models thus solved, we here adopt typical flows with short dynamical explosion time scales, $`\tau _{dyn}50\mathrm{ms}`$, by varying the entropy: ($`\tau _{dyn}`$ \[msec\], $`s/k`$)=(5, 355), (10, 480), (20, 750), (30, 1050), and (50, 1680), where we fixed the other parameters. These models produce almost the same final abundance of r-process elements (Sasaqui, Kajino and Mathews, 2005) and can well explain the abundance pattern of neutron-capture elements which were detected in one of the most r-process element enhanced metal-deficient halo star CS22892-052 (Sneden et al., 1996).
### 3.3 Neutrino Heating and Cutoff Time
Otsuki et al. (2000) showed that neutrino heating occurs most effectively at $`r30`$ km from the center of the collapsing core. We found in the same flow model analysis that it takes $`34`$ ms for the material blowing off the surface from proto-NS to reach 30 km. We discussed in sect. 2 that neutrino emission is followed until abruptly terminated by the BH formation with the appearance of apparent horizon in the first scenario. Finite time elapses before BH forms, and the neutrino luminosity is cutoff at $`12`$ s (Burrows, 1988). In the second scenario the cutoff time is delayed about 10 s for the formation of proto-NS, followed by possibly a phase transition of softening the EOS of core matter, and possibly later mass accretion onto BH. Since the neutrino cutoff time $`12`$ s or 10 s is larger than the heating time scale $`34`$ ms, we can assume that the hot bubble may form in the neutrino-driven wind of Type-II SNe that form BH as a remnant.
The temperature of all flows of the neutrino-driven wind used in the present study reaches to $`T_99`$ at $`t=34`$ ms. We adopt this temperature as the typical temperature at which neutrino heating is completed. The entropy of the hot bubble formed stays constant after this time. This justifies our assumption of constant entropy during the nucleosynthesis which follows. Our nucleosynthesis calculation starts from this initial temperature, and time zero refers to the time when the hot bubble reaches $`T_99`$ as displayed in Figure 1. It is to be noted that the time in this figure and all others we discuss later is not the time after core collapse or core bounce. Flows of the neutrino-driven wind successively blow off until the neutrino luminosity is cutoff at the time $`t=t_{cut}`$.
## 4 Nucleosynthesis Network
For the calculation of the r-process nucleosynthesis, we employ the reaction network used in Sasaqui et al. (2005a, b), which was developed from the original dynamical network code calculations described in Meyer et al. (1992), Woosley et al. (1994), Terasawa et al. (2001), and Otsuki et al. (2003).
The main feature of the network code developed (Sasaqui et al., 2005a, b) is the improvement in the light-mass nuclear reactions. It has already been discussed in literature (Woosley & Hoffman, 1992; Meyer et al., 1992; Woosley et al., 1994) that the $`\alpha `$($`\alpha `$n,$`\gamma `$)<sup>9</sup>Be($`\alpha `$,n)<sup>12</sup>C reaction sequence is particularly important in the earlier stage of the $`\alpha `$-process at high temperature and high density for the production of r-process seed nuclei. Terasawa et al. (2001) suggested that as long as the expansion timescales of the neutrino-driven winds are short, $`\tau _{dyn}10`$ms, a successful r-process occurs (Otsuki et al., 2000). Another nuclear reaction-flow path along the neutron-rich unstable nuclei may play as significant a role as the $`\alpha `$($`\alpha `$n,$`\gamma `$)<sup>9</sup>Be reaction does. Sasaqui et al. (2005a, b) found that this is always the case regardless of the flow models of neutrino-powered SN explosions. They quantitatively identified that the $`\alpha `$(t,$`\gamma `$)<sup>7</sup>Li reaction and the subsequent <sup>7</sup>Li(n,$`\gamma `$)<sup>8</sup>Li($`\alpha `$,n)<sup>11</sup>B reaction are the most critical reactions in addition to $`\alpha `$($`\alpha `$n,$`\gamma `$)<sup>9</sup>Be($`\alpha `$,n)<sup>12</sup>C that affect strongly the r-process nucleosynthesis particularly of the actinide elements, <sup>232</sup>Th and <sup>235,238</sup>U. They also found that these different nuclear reaction paths merge at <sup>14</sup>C, which is followed by neutron-capture flows on carbon, nitrogen and oxygen isotopes to manifest a new feature of “semi-waiting” point at the neutron-rich isotopes <sup>16</sup>C, <sup>18</sup>C, and <sup>24</sup>O. This feature of the “semi-waiting” point is the manifestation of a primary r-process so that the SN nucleosynthesis starts from the high entropy conditions on which the NSE favors neutrons, protons, and some amount of $`\alpha `$-particles as an initial composition and thereby the $`\alpha `$-capture, neutron-capture, and $`\beta `$-decay compete with one another in the light-mass neutron-rich nuclei. Sasaqui et al. (2005a, b) updated many nuclear reaction rates on light-mass neutron-rich nuclei with the help of recently accumulated new experimental data obtained by using radioactive nuclear beams (Nakamura et al., 1999; Sasaqui et al., 2005a, b) and references therein.
We also note that we calculate the nucleosynthesis sequence from the NSE, $`\alpha `$-process, $`\alpha `$-rich freeze-out, r-process, and subsequent beta-decay and alpha-decay, as explained in sect. 1 and in Figure 3, in a single network code rather than to split the calculation into two parts as was done in Woosley et al. (1994). This is important for our present study in looking the observable signatures of BH formation in elemental abundances from the r-process nucleosynthesis.
## 5 Result and Discussions
R-process nucleosynthesis in neutrino-driven wind of Type-II SN explosions has been studied by several authors (Woosley et al., 1994; Witti et al., 1994; Meyer, 1995; Qian & Woosley, 1996; Otsuki et al., 2000). The following two conditions prove to be important for a successful r-process:
$`\tau _{dyn}\tau _\nu ,`$ (17)
$`\tau _{dyn}\tau _{\alpha \alpha n},`$ (18)
where $`\tau _\nu 0.201\times \mathrm{L}_{\nu ,51}^1(\frac{ϵ_\nu }{\mathrm{MeV}})(\frac{\mathrm{r}}{100\mathrm{k}\mathrm{m}})^2(\frac{\sigma _\nu }{10^{41}\mathrm{cm}^2})^1`$s is the neutrino collision time scale (Qian et al., 1997), $`\tau _{\alpha \alpha n}[\rho ^2Y_\alpha ^2Y_n\sigma v_{\alpha \alpha n}N_A^2]^1`$ is the typical nuclear reaction time scale for <sup>4</sup>He$`(\alpha n,\gamma )^9`$Be which is the slowest one among all charged particle reaction paths, $`\mathrm{L}_{\nu ,51}`$ is the neutrino flux normalized in units of $`10^{51}`$ erg$`\mathrm{s}^1`$, $`ϵ_\nu `$ is the averaged electron-type neutrino energy $``$ 11 MeV, $`r`$ is the distance of r-process site from the center of the core, $`\sigma _\nu `$ is the averaged cross section over the neutrino energy spectrum, $`\rho `$ is the density of the mass element, and $`N_A`$ is Avogadro’s number. The former relation (13) describes the condition in which the neutron abundance is kept high when the neutrino process becomes ineffective (see discussion in the introduction), and the latter relation (14) describes the condition in which the neutron-to-seed ratio ($`Y_n/Y_{seed}`$) is high enough to realize a successful r-process of heavy elements from the seed nuclei. $`Y_{seed}`$ is defined as the total seed abundance, i.e. $`Y_{seed}=Y_A`$ for 70 $``$ A $``$ 120.
The time evolution of the temperature $`T_9`$, neutron separation energy $`S_n`$, and neutron and seed abundances $`Y_\mathrm{n}`$, $`Y_{\mathrm{seed}}`$, and their ratio $`Y_n/Y_{seed}`$ is shown in Figure 1 and Figure 2, separately. Here, $`S_n`$ means optimum neutron separation energy. This value is convenient to identify when the classical neutron-capture flow proceeds at low and almost constant value, $`S_n1`$ MeV, and when the freezeout of the n-capture process occurs. Assuming that the transition probabilities for (n, $`\gamma `$) and ($`\gamma `$,n) are equal to each other and nuclear reaction flow stays at a certain nucleus, $`S_n`$ is given by
$`S_n={\displaystyle \frac{T_9}{5.040}}\{34.075\mathrm{log}(Y_n\rho N_A)+{\displaystyle \frac{3}{2}}\mathrm{log}T_9\}.`$ (19)
When $`Y_n`$ is large in neutron-rich environment, $`S_n`$ is generally small. A drastic increase in $`S_n`$ gives us the information on neutron consumption at $``$ 1 s as shown in the bottom panel of Figure 1.
The calculated final abundance patterns are shown in Figure 3 for various values of neutrino cutoff time $`t_{cut}=0.001,0.005,0.1`$ s, and $`\mathrm{}`$. The last case is for no neutrino cutoff. The result is summarized as follows:
(1) When the neutrino cutoff occurs at $`t_{cut}=0.001`$ s under NSE condition (see Figure 1), the influence of the cutoff is largest as neutron abundance $`Y_n`$ is kept high (Figure 2) and the r-process becomes effective. This is because the neutrino interactions on neutrons occur for a short time before $`t_{cut}`$. As shown by the dashed line in Figure 3, the heavy r-process elements such as actinides are very effectively produced. However, in this case the cutoff time $`t_{cut}=1`$ ms is too short for shock wave to break out the iron core and SN might fail to explode due to the lack of neutrino heating.
(2) When the neutrino cutoff occurs at $`t_{cut}=0.005`$ s under the condition of efficient $`\alpha `$-process (see Figure 1), the effect of the neutrino cutoff is still large. The dotted line in Figure 3 shows that the r-process proceeds as effectively as in the first case, and the dotted and dashed lines are too close to each other to be separately read out. Both neutron abundance $`Y_n`$ and the neutron-to-seed ratio $`Y_n/Y_{seed}`$ in Figure 2 are almost the same as those in the first case.
(3) When the neutrino cutoff occurs at $`t_{cut}=0.1`$ s under the condition of efficient neutron-capture process (see Figure 1), the effect of the neutrino cutoff is small, and the final r-process abundance denoted by the dash-dotted line in Figure 3 is very close to the result of no neutrino cutoff (solid line). A notable difference from the previous two cases (1) and (2) is seen in $`Y_n`$ and $`Y_n/Y_{seed}`$ of Figure 2 for $`t_\mathrm{n}<t_{cut}`$.
(4) When we do not take account of the neutrino cutoff for $`t_{cut}=\mathrm{}`$, calculated <sup>232</sup>Th abundance is close to the observed lower limit from several metal-deficient halo stars as summarized in Table 3. Both $`Y_n`$ and neutron-to-seed ratio $`Y_n/Y_{seed}`$ in Figure 2 are close to those of the case (3).
The neutrino cutoff occurring in various nucleosynthesis stages shown in Figure 3 can thus make a remarkable effect on the actinide production. This is an important key feature to identify which SN remnant, BH or SN, forms in the gravitational core-collapse Type-II SN explosions.
We pay special attention to the behavior of the actinides (<sup>232</sup>Th, <sup>235,238</sup>U) because of their importance in the cosmochronology. Table 1 and Figure 4 show the ratio of <sup>232</sup>Th/(<sup>151</sup>Eu+<sup>153</sup>Eu). This is a useful quantity because astronomical observations of metal-deficient halo stars cannot provide each isotopic abundance but just this ratio. We find in Figure 3 that the early cutoff time makes this ratio large. We also find that there appear two outcomes separated by a narrow transitional neutrino cutoff time between 0.01 s and 0.1 s. The “BH outcome” shows a very high abundance ratio which is thought to arise from the events of BH formation, while the “NS outcome” shows a low abundance ratio which may arise from the conditions that do not lead to BH formation. Also, as discussed in sect. 1, the effect of the neutrino cutoff by BH formation is mainly due to the change of neutrons into protons by the weak reaction process (4). So, the drastic change of the final abundance pattern occurs if the neutrino cutoff occurs at a time right after the end of the $`\alpha `$-process and the beginning of the neutron-capture process, i.e. $`t_nt`$. This is because drastic environmental change occurs after the $`\alpha `$-rich freezeout of making seed elements at $`t_nt`$ and only abundant neutrons are easily affected by the weak process (4).
This profile can be fit by the following function.
$`z=z_t\mathrm{}+{\displaystyle \frac{z_1}{1+(\frac{t}{t_0})^\alpha }},`$ (20)
where $`z^{232}`$Th/(<sup>151</sup>Eu+<sup>153</sup>Eu), $`tt_{cut}`$, $`z_t\mathrm{}`$ means the abundance ratio of no neutrino cutoff, $`z_1+z_t\mathrm{}(z_0)`$ means the abundance ratio with no neutrinos, $`t_0`$ means the start of the drastic change of the ratio $`z`$, and $`\alpha `$ is a constant value. The quantity $`\alpha `$ is expected to be model independent and the other quantities are model dependent. In case of Figure 4, for example, $`z_t\mathrm{}`$ is 0.284, $`z_1`$ is 5.20, $`t_0`$ is 0.0199 s, and $`\alpha `$ is 3.00 (Figure 4) in Otsuki model (Otsuki et al., 2000).
We repeated this calculation using the exponential flow models which are characterized by the different dynamical expansion time scales $`\tau _{dyn}=5,10,20,30`$, and 50 ms, as defined in sect. 3.2. The calculated results are shown in Table 2 and Figure 5 in normalized form of Eq. (16)
$`\xi ={\displaystyle \frac{zz_t\mathrm{}}{z_1}}={\displaystyle \frac{1}{1+(t/t_0)^\alpha }}.`$ (21)
Here, we set $`\alpha `$, the parameter which determines the shape of fitting function, to be 3, the same as that of Figure 4. The parameter $`t_0`$ in Eq. (16), which corresponds to the time around which a drastic change of the final abundance occurs, is proportional to the dynamical timescale $`\tau _{dyn}`$. In the models with a small dynamical timescale, the temperature drops rapidly and the charged particle reactions are suppressed soon. The freeze out of the $`\alpha `$-process occurs early, so the neutrino cutoff effect leading to a drastic change in the abundance ratio $`z`$ occurs early. On the other hand, in the models with a large dynamical timescale, the temperature drops slowly and the charged particle reactions continue to take place. The time of the freezeout of the $`\alpha `$-process becomes later than that in the models with the small dynamical timescale. The drastic change occurs mostly for this reason.
Beacom et al. (2001) proposed that the direct possible signature of BH formation could be the observation of sharp cutoff in the neutrino signal. This may not be a unique good method because there is increasing evidence that the BH formation must be rare at the present epoch for ordinary core-collapse SN scenario, given the low frequency of nearby SNe. Strigari et al. (2003) have shown that direct measurements of the evolution of core-collapse SN rate, which includes the SN rate for BH formation, are consistent with the predictions based on a variety of star formation indicators. From the view point of Galactic chemical evolution, however, different physics may operate (Heger et al., 2003) in the SN explosions in various environments with very different metalicities from the early Galaxy to the present epoch. Neither the neutrino detection studies nor the star formation rate studies can address the possibility that BH formation was a more frequent and dynamically important process in the very early Galaxy of active star formation epoch.
The nucleosynthesis signal proposed in this paper reveals the promising possibility that the effect of neutrino cutoff associated with BH formation would manifest itself in the fossil record of the produced r-process yields. Following remarkable advances in spectroscopic observations, significant information about the abundance of the r-process elements in metal-deficient halo stars has recently been accumulated. These stars are presumed to be second generation stars which were born in the early Galaxy and their elemental abundances were ideally affected by a few SN episodes of the first generation massive population III stars.
In Figure 3 we show the observed abundance ratios <sup>232</sup>Th/(<sup>151</sup>Eu+<sup>153</sup>Eu) and
(<sup>235</sup>U+<sup>238</sup>U)/(<sup>151</sup>Eu+<sup>153</sup>Eu) of several metal-deficient halo stars tabulated in Table 3. It is interesting that the typical predicted abundances are near the low end of the observed yield in this Figure. This might possibly indicate that those actinides were produced in the first generation SNe associated with remnant NS formation or indicate that only a partial contribution from the SN products associated with BH formation is admixed in the observed r-process elements. The same result is more clearly shown in Figure 4 and 5 which display a band of observed abundance ratios near the “NS outcome” rather than the “BH outcome”.
## 6 Summary and Future Outlook
We investigated the r-process nucleosynthesis in the SNe where BH could form. We found that the r-process abundances could change significantly by neutrino cutoff at the BH formation. This will be one of the predictions about BH formation if metal-deficient halo dwarfs which have such a specific abundance pattern can be found. There appear no observational signatures at the moment indicating that the abundance pattern is made in the events where the BH can be formed. Future observations of more halo stars exhibiting enhanced r-process elements are highly desirable.
We assumed steady state flow of the neutrino-driven wind in our present study. Actually, the flow model must be effected by the BH formation. Therefore reliable models with realistic numerical simulations of general relativistic hydrodynamics are needed. However, the dynamics would be more complex since a very massive star, which evolves to a standard SN, should collapse by accreting materials from the accretion disk and eject part of them in a jet-like explosion. We need to know more details about the dynamics where the BH can be formed. We also need to understand the behavior of the neutrino-driven wind in the environments where the neutrino luminosity is cutoff.
Massive stars associated with BH formation are predicted to occur in the early stage of the Galactic evolution so that they could have ejected nucleosynthesis products which are very different from those ejected from ordinary SNe that leave NS as a remnant. In this article we proposed that actinides could show a remarkable difference for the different neutrino cutoff effects in SNe leaving BH or NS. Similar differences would be found also in the elemental abundances of <sup>7</sup>Li and <sup>11</sup>B (Yoshida et al., 2004, 2005) which are produced in the neutrino processes in the outer layers of core-collapse SNe. Such productions could impact the Galactic chemical evolution, a point worth investigating. The ultra metal-deficient stars which have such an abundance pattern might be found in the future using more sophisticated observational techniques.
This work has been supported in part by Grants-in-Aid for Scientific Research (13640313, 17540275) and is for Specially Promoted Research (13002001) of the Ministry of Education, Science, Sports and Culture of Japan, and The Mitsubishi Foundation. This work has also been supported in part by the U.S. National Science Foundation Grant No. PHY-0244384 and by the University of Wisconsin Research Committee with funds granted by the Wisconsin Alumni Research Foundation. A.B.B. gratefully acknowledges the 21st Century for Center of Excellence Program “Exploring New Science by Bridging Particle-Matter Hierarchy” at Tohoku University for financial support and thanks the Nuclear Theory Group at Tohoku University for their hospitality. |
warning/0506/hep-th0506074.html | ar5iv | text | # Untitled Document
DCPT-05/29
A quasi-particle description of the $`(3,p)`$ models
P. Jacob and P. Mathieu
Department of Mathematical Sciences, University of Durham, Durham, DH1 3LE, UK
and
Département de physique,Université Laval, Québec, Canada G1K 7P4
(patrick.jacob@durham.ac.uk, pmathieu@phy.ulaval.ca)
Abstract
The $`(3,p)`$ minimal models are reconsidered from the point of view of the extended algebra whose generators are the energy-momentum tensor and the primary field $`\varphi _{2,1}`$ of dimension $`\left(p2\right)/4`$. Within this framework, we provide a quasi-particle description of these models, in which all states are expressed solely in terms of the $`\varphi _{2,1}`$-modes. More precisely, we show that all the states can be written in terms of $`\varphi _{2,1}`$-type highest-weight states and their $`\varphi _{2,1}`$-descendants. We further demonstrate that the conformal dimension of these highest-weight states can be calculated from the $`\varphi _{2,1}`$ commutation relations, the highest-weight conditions and associativity. For the simplest models $`\left(p=5,\mathrm{\hspace{0.17em}7}\right)`$, the full spectrum is explicitly reconstructed along these lines. For $`p`$ odd, the commutation relations between the $`\varphi _{2,1}`$ modes take the form of infinite sums, i.e., of generalized commutation relations akin to parafermionic models. In that case, an unexpected operator, generalizing the Witten index, is unravelled in the OPE of $`\varphi _{2,1}`$ with itself. A quasi-particle basis formulated in terms of the sole $`\varphi _{1,2}`$ modes is studied for all allowed values of $`p`$. We argue that it is governed by jagged-type partitions further subject a difference 2 condition at distance 2. We demonstrate the correctness of this basis by constructing its generating function, from which the proper fermionic expression of the combination of the Virasoro irreducible characters $`\chi _{1,s}`$ and $`\chi _{1,ps}`$ (for $`1s\left[p/3\right]+1`$) are recovered. As an aside, a practical technique for implementing associativity at the level of mode computations is presented, together with a general discussion of the relation between associativity and the Jacobi identities.
05/05
1. Introduction
1.1. Quasi-particle description of the minimal models: extended algebra vs spinon-type formulation
Fermionic-type character expressions are known for all Virasoro minimal models.<sup>1</sup> A brief review of the origin of fermionic characters in conformal field theory, together with an extended list of references (focusing mainly on the early works), is presented in the introduction of . Note that the first fermionic character formulas appeared in mathematics and its interpretation in terms of an exclusion principle seems to go back to . Another early reference on fermionic formulas is . This subject has somewhat exploded with where various fermionic-sum formulas were proposed and interpreted in terms of a generalized exclusion principle. More formulas pertaining to the Virasoro minimal models were conjectured and proved in and . The works of the last reference rely heavily on a special truncation of the spin chain XXZ spectrum. A different method for proving the Virasoro fermionic formulas has been developed in , using the connection with the RSOS models . The recent work contains exhaustive results and many references to other works. A different approach to the construction of fermionic formulas has been initiated in and vigorously extended recently by Feigin and collaborators. However, most of these Virasoro characters, as well as the underlying basis of states, have not been explained within a conformal field theoretical set up. To look for such an intrinsic understanding of the fermionic characters is a fundamental quest. So far, only the minimal models $`(2,p)`$ have been successfully addressed from that perspective . But is there any natural lines of attack for handling this question? Two potential avenues can readily been identified: a reinterpretation in terms of an extended algebra or a spinon-type reformulation.
All minimal models have a hidden extended conformal symmetry. The extension is obtained by adding one generator to the usual energy-momentum tensor. This extra generator is $`\varphi _{1,p1}=\varphi _{p^{}1,1}`$,<sup>2</sup> The field $`\varphi _{p^{}1,1}`$ has dimension $`h_{p^{}1,1}=(p2)(p^{}2)/4`$, which is integer when either $`p`$ or $`p^{}`$ is of the form $`2+4m`$. In that case, this extended algebra is at the roots of the corresponding A-D type block-diagonal modular invariants (for this interpretation, see for instance ). Notice that even for this case, the minimal models have not been reformulated in terms of the representation theory of this extended algebra. But we stress that such a reformulation does not requires $`h_{p^{}1,1}`$ to be integer (a point that is plainly illustrated in this work). with fusion rule
$$\varphi _{p^{}1,1}\times \varphi _{p^{}1,1}=\varphi _{1,1}=I.$$
$`(1.1)`$
Reformulating the minimal models from the perspective of this two-generator extended algebra has the obvious disadvantage that the defining extended algebra differs from one model to the other. But on the other hand, that the algebraic formalism is tailor-made for each model strongly suggests the existence of an underlying basis, generated by the $`\varphi _{p^{}1,1}`$ modes, that would be free of singular vectors, i.e., a quasi-particle basis.
An alternative formulation might be considered, in which the fundamental spectrum generating field is either $`\varphi _{1,2}`$ or $`\varphi _{2,1}`$. In such a formulation, the extended algebra (incorporating one of $`\varphi _{1,2}`$ and $`\varphi _{2,1}`$) would generically be ‘non-abelian’ (that is, multi-channels), e.g., since
$$\varphi _{2,1}\times \varphi _{2,1}=\varphi _{1,1}+\varphi _{3,1},$$
$`(1.2)`$
with the two fields on the right-hand side having dimensions that do not differ by an integer.<sup>3</sup> They differ by an integer when $`p^{}=2`$ but in that case $`\varphi _{3,1}`$ is a descendant of $`\varphi _{1,1}`$ . In that framework, the added algebra generators would be the set of fields $`\{\varphi _{r,1}\}`$ in which $`\varphi _{2,1}`$ plays the role of a basic generator. This is much like the parafermionic algebra which is spanned by the parafermionic fields $`\psi _n`$, with $`0nk1`$; the $`\psi _n`$’s can all be obtained from the multiple product of $`\psi _1`$ with itself. However, the parafermionic algebra is single-channel (that is, $`\psi _n\times \psi _m=\psi _{n+m}`$), in contradistinction with the above $`\{\varphi _{r,1}\}`$ algebra. This yet-to-be-defined $`\varphi _{2,1}`$-type reformulation of the minimal models is actually closer to the spinon description of the $`\widehat{su}(2)_k`$ WZW models (see for $`k=1`$ and for $`k>1`$). In the WZW context, the spinon field refers to the primary doublet of spin $`j=1/2`$, denoted $`\varphi _{1/2}`$ and the analog of the above fusion rule is $`\varphi _{1/2}\times \varphi _{1/2}=\varphi _0+\varphi _1`$ ($`\varphi _1`$ being absent for $`k=1`$). Because of this analogy, this approach will be referred to as a spinon-type reformulation.
Extended algebras associated to a set of OPE including multi-channel ones of the sort (1.2) appear rather difficult to analyze, however. But that the field $`\varphi _{1,2}`$ could be regarded as a primary field in a formulation based on an extended algebra containing $`\varphi _{2,1}`$ (or the inverse) would fit nicely the following suggestive formula that would necessarily results from such a construction
$$\frac{2h_{1,2}h_{2,1}}{c_{p^{},p}}=\frac{1}{8}.$$
$`(1.3)`$
This expression is valid for all minimal models $`(p^{},p)`$, with
$$h_{r,s}=\frac{(rpsp^{})^2(pp^{})^2}{4pp^{}}\mathrm{and}c_{p^{},p}=1\frac{6(pp^{})^2}{pp^{}},$$
$`(1.4)`$
($`1rp^{}1`$ and $`1sp1`$, with $`p^{}<p`$).
Whether this spinon-type approach can be worked out in general remains to be seen. In this work, we consider the special case for which this second method reduces to the first one, namely the $`(3,p)`$ minimal models. When $`p^{}=3`$, the fusion of $`\varphi _{2,1}`$ with itself yields the identity:
$$\varphi _{2,1}\times \varphi _{2,1}=I.$$
$`(1.5)`$
1.2. Generalities on the structure of the $`\varphi `$-algebra
Our aim here is thus to describe the $`(3,p)`$ models in terms of the extended algebra associated to the OPEs
$$\begin{array}{cc}\hfill \varphi (z)\varphi (w)=& \frac{1}{(zw)^{2h}}\left[I+(zw)^2\frac{2h}{c}T(w)+\mathrm{}\right]𝒮,\hfill \\ \hfill T(z)\varphi (w)=& \frac{h\varphi (w)}{(zw)^2}+\frac{\varphi (w)}{(zw)}+\mathrm{}\hfill \\ \hfill T(z)T(w)=& \frac{c_{3,p}/2}{(zw)^4}+\frac{2T(w)}{(zw)^2}+\frac{T(w)}{(zw)}+\mathrm{}\hfill \end{array}$$
$`(1.6)`$
with
$$\varphi \varphi _{2,1}\mathrm{and}hh_{1,2}=\frac{p2}{4},$$
$`(1.7)`$
instead of through the irreducible representations of the Virasoro algebra. In the first OPE, we have included an operator $`𝒮`$ which is enforced (see below) by the mutual locality of $`\varphi `$ for the cases where $`h\mathrm{}_+/2`$, that is, for $`p`$ odd; $`𝒮`$ anticommutes with $`\varphi `$ and commutes with $`T`$. That this algebra is associative , at least for the value central charge $`c_{3,p}`$, is guaranteed by the associativity of the operator algebra of minimal models .
Note that, as for the WZW models or the parafermionic models , the Virasoro algebra lives in the ‘enveloping algebra’ of $`\varphi `$. This is well-known for the free-fermionic formulation of the Ising model, which corresponds to the $`p=4`$ case of the above construction (i.e., $`\psi =\varphi _{1,3}=\varphi _{2,1}`$). That $`T`$ might not occur in a (pole-type) singular term in the OPE $`\varphi (z)\varphi (w)`$ (which is the case for $`p<6`$) is no obstruction to the fact that $`T`$ can always be written as a bilinear in the field $`\varphi `$. In that sense, the ‘$`\varphi `$-algebra’ is essentially defined by the first OPE in (1.6) and it will be understood as such.
We stress that, in writing (1.6), we do not consider the generic form of the OPE $`\varphi (z)\varphi (w)`$ . In fact, we have imposed an implicit restriction on those terms that can appear on the right-hand side of (1.6): these are solely the descendants of the Virasoro identity. In particular, no bilinear composite of $`\varphi `$ does occur. This condition is restrictive; it actually fixes $`c`$ to finitely many values (and in particular, $`c`$ cannot be free). For instance, when $`\varphi `$ has dimension $`h=3/2`$, there are only two solutions to the associativity conditions: $`c=7/10`$ and $`c=21/4`$.
The structure of the $`\varphi `$-algebra (1.6) depends crucially upon the parity of $`p`$. For $`p`$ even, $`2h`$ is integer so that $`𝒮=I`$. There are then only integer powers of $`zw`$ on the right-hand side of (1.6). This results into ordinary (anti)commutation relations. The basic data for the first few models with $`p`$ even are:
| | $`(3,8):`$ | $`h=3/2,`$ | $`c=21/4:𝒮(2,8),`$ | $`(3,16):`$ | $`h=7/2,`$ | $`c=161/8,`$ |
| --- | --- | --- | --- | --- | --- | --- |
| | $`(3,10):`$ | $`h=2,`$ | $`c=44/5:[(2,5)]^2,`$ | $`(3,20):`$ | $`h=9/2,`$ | $`c=279/10,`$ |
| | $`(3,14):`$ | $`h=3,`$ | $`c=114/7:W_3(3,7),`$ | $`(3,22):`$ | $`h=5,`$ | $`c=350/11.`$ |
$`(1.8)`$
$`𝒮`$ stands for a superconformal minimal model. That associativity could fix the central charge in spite of the fact that these relations might satisfy the Jacobi identities for generic values of the central charge, as in the first three cases , is a point that is clarified in the appendix. The solution displayed in the last three cases is among those few (3 for $`h=7/2`$ and 5 for the other two cases) found in an investigation of low-dimensional two-generators extended algebras .
On the other hand, for $`p`$ odd, $`h=\pm 1/4`$ mod 1. That case presents at once a severe complication in that the commutation relations resulting from (1.6) take the form of infinite sums, as for parafermionic models (cf. section 2). Our analysis of the first two models in this class, $`h=3/4`$ and $`h=5/4`$, shows that in these cases $`c`$ is uniquely determined and it agrees with $`c_{3,5}`$ and $`c_{3,7}`$ respectively. But that two solutions are found for the next case ($`h=9/4`$) suggests that this uniqueness is not generic.
As previously mentioned, the analysis of the $`p`$ odd case also reveals that the first OPE in (1.6) requires the introduction of the operator $`𝒮`$, anticommuting with $`\varphi `$. This operator can be viewed as a generalization of the Witten operator $`(1)^F`$ which anticommutes with fermions . Here we have $`𝒮=(1)^p`$ with
$$\varphi \times \varphi =(1)^pI\mathrm{with}(1)^p\varphi =(1)^p\varphi (1)^p.$$
$`(1.9)`$
This is quite similar to the phase factor introduced in the $`\widehat{su}(2)_1`$ commutation relations for the spinon fields (for which $`h=1/4`$) in (cf. eq. (3) there). The existence of this operator is definitely forced by associativity. Note that it does not appear in the OPE of the physical field $`\mathrm{\Phi }(z,\overline{z})=\varphi (z)\varphi (\overline{z})`$ since it squares to 1. Since $`T`$ is bilinear in $`\varphi `$, it commutes with $`𝒮`$.
1.3. The spinon-type reformulation of the $`(3,p)`$ models: setting the problem
What does a $`\varphi `$-algebra reformulation of the $`(3,p)`$ models should amount to? It should certainly allow us to fix completely the spectrum of the model for a given $`p`$, that is, to determine the highest-weight states and their conformal dimension. The highest-weight states turn out to be completely characterized by an integer $`\mathrm{}`$ such that $`0\mathrm{}(p2)/2`$. The highest-weight state conditions are formulated directly in terms of the $`\varphi `$-modes and they read
$$\varphi _{hn+\frac{\mathrm{}}{2}}|\sigma _{\mathrm{}}=0n>0.$$
$`(1.10)`$
We then observe that the conformal dimension of some highest-weight states $`|\sigma _{\mathrm{}}`$ are directly obtained from the commutation relations deduced from (1.6) and the highest-weight conditions (1.10). However, it is generally necessary to invoke associativity to fix the full spectrum. This analysis is performed in section 2 for the two simplest $`(3,p)`$ models apart from the Ising case, namely $`(3,5)`$ and $`(3,7)`$. The obtained conformal dimensions confirm the following identification with the Virasoro highest-weight states: $`|\varphi _{1,s}=|\sigma _{s1}`$ for $`1s<p/2`$.
A second required ingredient is a complete characterization of the descendant states in terms of the action of the $`\varphi `$-modes. Note in that regard that a single $`\varphi `$-module generically decompose into the direct sum of two Virasoso modules. Each Virasoro component will be singled out from the subset of descendant states containing an even or odd number of $`\varphi `$-modes acting on the highest-weight state. This feature is well known for those $`(3,p)`$ models that have previously been formulated from that perspective, namely the $`(3,4),(3,5)`$ and $`(3,8)`$ models, for which $`\varphi `$ is respectively the free fermion, the fundamental graded parafermion and the superpartner of the energy-momentum tensor. In the present case, the Virasoro highest-weight states $`|\varphi _{1,s}=|\sigma _{s1}`$ and $`|\varphi _{1,ps}`$ ($`1s<p/2)`$ are combined into a single $`\varphi `$-module since $`|\varphi _{1,ps}`$ is a descendant of $`|\varphi _{1,s}`$:
$$|\varphi _{1,ps}=\varphi _{h+(s1)/2}|\varphi _{1,s}.$$
$`(1.11)`$
This matches the relation
$$h_{1,ps}h_{1,s}=h\frac{(s1)}{2}.$$
$`(1.12)`$
We stress that we invoke here only very general aspects of a yet-to-be-defined representation theory of the $`\varphi `$-algebra. Actually, we make use the notion of a highest-weight state, that of $`\varphi `$-lowering operators and associativity.
Ideally, we would also like to look for a complete description of the irreducible modules, i.e., a basis of states. There are two natural possibilities for such a basis, one formulated in terms of the combination of the Virasoro modes and the $`\varphi `$-modes or one formulated solely in terms of the $`\varphi `$-modes. In this work, we focus on the second possibility. In that case, the basis is expected to be of a quasi-particle type. A quasi-particle basis entails the construction of the Hilbert space by the action of the quasi-particle creation operators subject to a restriction rule, i.e., a filling process controlled by an exclusion principle. The obtention of this basis is our main result.
1.4. Quasi-particle basis of states
In the quest for a quasi-particle basis, we were guided by our previous construction of the quasi-particle basis for the graded parafermions in terms of jagged partitions. Such partitions differ from standard partitions, for which parts are non-increasing from left to right, in that a possible increase between parts at short distance is allowed.
Manipulations with the generalized commutation relations derived from (1.6) together with simple considerations on the structure of the highest-weight modules lead us to infer the general form of a candidate basis (section 3). It takes the following form. The highest-weight modules $`|\sigma _{\mathrm{}}`$ are described by the successive action of the lowering $`\varphi `$-modes subject to specific constraints. In the $`N`$-particle sector, with strings of lowering modes written in the form
$$\varphi _{h+\frac{\mathrm{}}{2}+\frac{(N1)}{2}n_1}\varphi _{h+\frac{\mathrm{}}{2}+\frac{(N2)}{2}n_2}\mathrm{}\varphi _{h+\frac{\mathrm{}}{2}+\frac{1}{2}n_{N1}}\varphi _{h+\frac{\mathrm{}}{2}n_N}|\sigma _{\mathrm{}},$$
$`(1.13)`$
these constraints are:
| | $`n_in_{i+1}r+1`$ | $`n_in_{i+2}+2`$ |
| --- | --- | --- |
| | $`n_{N1}\mathrm{}r`$ | $`n_N0,`$ |
$`(1.14)`$
where
$$2r=p5.$$
$`(1.15)`$
The $`n_i`$’s are integers for $`p`$ odd and alternate between integer and half-integer values when $`p`$ is even (from right to left). The condition $`0\mathrm{}k`$ is linked to the boundary condition on $`n_{N1}`$ that appears to be ‘complete’ (i.e., to represent the full set of conditions that singles out the different modules) only in these cases.
This quasi-particle basis of states was known in at least three cases: $`p=4,5`$ and $`8`$. In each case, it reduces to (1.14). For $`p=4`$, the quasi-particle basis is that of a free fermion:
$$b_{s_1}\mathrm{}b_{s_N}|0\mathrm{with}s_is_{i+1}+1.$$
$`(1.16)`$
To rephrase this in terms of our previous notation, we set
$$s_i=n_i+\frac{1}{2}\frac{(Ni)}{2}n_in_{i+1}+\frac{3}{2},$$
$`(1.17)`$
which is indeed the first condition in (1.14) when $`r=1/2`$. In this special case, the second condition is a consequence of the first one.
For $`p=5`$, as mentioned previously, the model is the simplest example of a graded $`\mathrm{}_k`$ parafermion , corresponding to the value $`k=1`$. The quasi-particle basis in that special case, when reformulated in terms of the $`n_i`$’s reads (see , end of section 5):
$$n_in_{i+1}+1.$$
$`(1.18)`$
This is again the first condition in (1.14) for $`r=0`$; here again it implies the second condition of (1.14).
For $`p=8`$, it has been pointed out that $`(3,8)𝒮(2,8)`$. Now, we have recently obtained the quasi-particle basis of superconformal models $`𝒮(2,4\kappa )`$ in . It is expressed solely in terms of $`G`$ modes ($`G`$ being the superpartner of $`T`$) and for $`\kappa =4`$, it takes the form
$$G_{s_1}\mathrm{}G_{s_N}|0\mathrm{with}s_is_{i+1}1\mathrm{and}s_is_{i+2}+1,$$
$`(1.19)`$
with all $`s_i`$ half-integers (cf. eqs (13) and (17)). The relation between $`s_i`$ and $`n_i`$ being $`s_i=n_i+3/2(Ni)/2`$, the above conditions translate into
$$n_in_{i+1}\frac{1}{2}\mathrm{and}n_in_{i+2}+2.$$
$`(1.20)`$
We again recover (1.14) for $`r=3/2`$.
1.5. Fermionic characters
The complete module of $`|\sigma _{\mathrm{}}`$ is obtained by summing over all these states (1.13) satisfying (1.14) and all values of $`N`$. Granting the correctness of this basis, one can then enumerate states in highest-weight modules (section 4). This leads to a fermionic expression of the (normalized) character that takes the simple form
$$\widehat{\chi }_{\mathrm{}}(q)=\underset{m_1,m_2,\mathrm{}m_k0}{}\frac{q^{𝐦B𝐦+C𝐦}}{(q)_{m_1}\mathrm{}(q)_{m_k}},$$
$`(1.21)`$
where the matrices $`B`$ and $`C`$ are defined in section 4. In terms of the Virasoro characters, $`\widehat{\chi }_{\mathrm{}}`$ decomposes as follows.:
$$\widehat{\chi }_{\mathrm{}}(q)=q^{h_{1,s}+c/24}\left[\chi _{1,s}^{\mathrm{Vir}}(q)+q^{h_{1,ps}h_{1,s}}\chi _{1,ps}^{\mathrm{Vir}}(q)\right](\mathrm{}=s\mathrm{1\hspace{0.33em}0}\mathrm{}[p/3]).$$
$`(1.22)`$
We recover in this way the fermionic sums given in , where these expressions where first conjectured. The expression of some of these characters can also be found in (see also the second reference of ). Their derivation from the general expressions in is presented in . This equivalence confirms the correctness of the basis (1.13)-(1.14).
Obtaining the fermionic character amounts to finding the generating function for all the states (1.13)-(1.14). But finding such generating functions is in general a difficult problem. In the present case, we modify the characterization of our states in order to make use of a related generating function derived in . Given that we read off our generating function (up to boundary terms) from and that this latter article explicitly deals with an algebra related to the $`(3,p)`$ models, it is appropriate to clarify the relation between this work and the present one. The authors of construct a vertex operator algebra out of the product of the $`(3,p)`$ model and a free boson. The algebra is generated by the two local fields: $`a(z)=V_1\varphi _{2,1}`$ and $`a^{}(z)=V_1\varphi _{2,1}`$, where $`V_1`$ and $`V_1`$ are vertex operators with dimension such that $`a,a^{}`$ have respective dimension $`1`$ and $`p3`$. The monomial basis underlying this vertex operator algebra is spanned by the modes $`a_\lambda =(a_{\lambda _1},\mathrm{},a_{\lambda _N})`$, with the $`\lambda _i`$’s subject to $`\lambda _i\lambda _{i+2}+2r`$. The origin of the exclusion here is rooted in polynomial relations of the type $`a^2^na=0`$ for $`0np3`$ (plus an additional one). In view of our results, by stripping off the contribution of the free boson, we end up with the jagged-type basis (1.14). In preparing the revised version of this work, we became aware of where the translation of these results to the $`(3,p)`$ models is performed and the resulting basis (cf. Lemma 5.5 there) agrees perfectly with ours for $`\mathrm{}=0`$.
2. The $`(3,p)`$ algebra for $`p`$ odd
2.1. Generalized commutation relations
Let us first justify the necessity of the operator $`𝒮`$ for $`p`$ odd by invoking associativity. Consider a correlator of the form $`\varphi (z_1)\varphi (z_2)\varphi (z_3)\mathrm{}`$. The first OPE in (1.6) shows that moving $`\varphi (z_2)`$ and then $`\varphi (z_3)`$ in front of $`\varphi (z_1)`$ induces a negative phase (since $`2h`$ is half-integer for $`p`$ odd):
$$\varphi (z_1)\varphi (z_2)\varphi (z_3)\mathrm{}=\varphi (z_2)\varphi (z_3)\varphi (z_1)\mathrm{}\frac{1}{z_{23}^{2h}}\left(𝒮+\mathrm{}\right)\varphi (z_1)\mathrm{},$$
$`(2.1)`$
which is to be compared with
$$\varphi (z_1)\varphi (z_2)\varphi (z_3)\mathrm{}\frac{1}{z_{23}^{2h}}\varphi (z_1)(𝒮+\mathrm{})\mathrm{}.$$
$`(2.2)`$
The compatibility of these expressions forces
$$𝒮\varphi (z)=\varphi (z)𝒮,$$
$`(2.3)`$
which captures the whole effect of $`𝒮`$ and allows for the identification with $`(1)^p`$ previously displayed in (1.9). For the $`(3,5)`$ case, this operator also appears in the description of the model as a graded parafermionic theory. This is briefly reviewed in the following subsection.
Consider now the commutation relations associated to (1.6). The mode decomposition of $`\varphi `$ acting on a state of ‘charge’ $`\mathrm{}`$ (or, in the sector specified by the integer $`\mathrm{}`$) is:
$$\varphi (z)=\underset{n\mathrm{}}{}z^{n\frac{\mathrm{}}{2}}\varphi _{h+\frac{\mathrm{}}{2}+n}.$$
$`(2.4)`$
The field $`\varphi `$ itself has charge 1. A first useful commutation relation is obtained by evaluating the following integral:
$$\frac{1}{(2\pi i)^2}𝑑w𝑑z\varphi (z)\varphi (w)z^{\frac{\mathrm{}}{2}+n}w^{\frac{\mathrm{}}{2}+m1}(zw)^{2h1},$$
$`(2.5)`$
in two different ways . That yields
$$\underset{t=0}{\overset{\mathrm{}}{}}C_{2h1}^{(t)}[\varphi _{\frac{\mathrm{}}{2}+nt+h}\varphi _{\frac{\mathrm{}}{2}+m+th}+\varphi _{\frac{\mathrm{}}{2}+m1t+h}\varphi _{\frac{\mathrm{}}{2}+n+1+th}]=𝒮\delta _{n+m+\mathrm{},0},$$
$`(2.6)`$
where
$$C_u^{(t)}=\frac{\mathrm{\Gamma }(tu)}{t!\mathrm{\Gamma }(u)}.$$
$`(2.7)`$
If we replace
$$w^{\frac{\mathrm{}}{2}+m1}(zw)^{2h1}w^{\frac{\mathrm{}}{2}+m+1}(zw)^{2h3}$$
$`(2.8)`$
in (2.5), in order to pick up the contribution of $`T`$ (the change in the power of $`w`$ being purely conventional), we get instead
$$\begin{array}{cc}\hfill \underset{t=0}{\overset{\mathrm{}}{}}C_{2h3}^{(t)}& [\varphi _{\frac{\mathrm{}}{2}+n2t+h}\varphi _{\frac{\mathrm{}}{2}+m+th+2}+\varphi _{\frac{\mathrm{}}{2}+m1t+h}\varphi _{\frac{\mathrm{}}{2}+n+1+th}]\hfill \\ & =\frac{1}{2}\left(\frac{\mathrm{}}{2}+n\right)\left(\frac{\mathrm{}}{2}+n1\right)𝒮\delta _{n+m+\mathrm{},0}+\frac{2h}{c}L_{n+m+\mathrm{}}𝒮.\hfill \end{array}$$
$`(2.9)`$
In the following, we will refer to the above two generalized commutation relations as follows:
$$\begin{array}{cc}& (2.6)\mathrm{I}_{n,m,\mathrm{}}\hfill \\ & (2.9)\mathrm{II}_{n,m,\mathrm{}}.\hfill \end{array}$$
$`(2.10)`$
Let us first test the relative sign for the two terms in the infinite sum (2.9)<sup>4</sup> Note that this sign is correlated to that in (2.6). This computation verifies the ‘bosonic’ nature of the field $`\varphi `$ within the present framework, i.e., that the interchange of the two fields in (2.5) does not generate a minus sign. by acting with both sides of (2.9) on the vacuum state $`|0`$ (so that $`\mathrm{}=0`$) which is such that
$$\varphi _{h+n}|0=0n>0.$$
$`(2.11)`$
With $`n=2`$ and $`m=2`$, only one term contributes from the first sum and none from the second sum, so that
$$\varphi _h\varphi _h|0=𝒮|0.$$
$`(2.12)`$
Taking instead $`n=m=0`$, only one term of the second sum contributes and we get the same result, confirming thus the positive relative sign.
We have just stated in (2.11) the highest-weight condition pertaining to the vacuum state. Let us denote by $`|\sigma _{\mathrm{}}`$ the highest-weight state in the sector labeled by $`\mathrm{}`$. Its highest-weight state characterization is
$$\varphi _{hn+\frac{\mathrm{}}{2}}|\sigma _{\mathrm{}}=0n>0.$$
$`(2.13)`$
Note that in order for the dimension of the first descendant $`\varphi _{h+\frac{\mathrm{}}{2}}|\sigma _{\mathrm{}}`$ to be non-negative, we require
$$0\mathrm{}\frac{p2}{2}.$$
$`(2.14)`$
This bound will be assumed to hold from now on. The action of $`𝒮`$ on a highest-weight state $`|\sigma _{\mathrm{}}`$ is normalized as
$$𝒮|\sigma _{\mathrm{}}=|\sigma _{\mathrm{}}.$$
$`(2.15)`$
A somewhat remarkable feature of the $`\varphi `$-algebra for $`p`$ odd is that the dimension of the highest-weight state $`|\sigma _1`$ follows directly from (2.9) and (2.13), exactly as for all parafermionic highest-weight states . (This, of course, is also true for $`\sigma _0=I`$.) Applying both sides of (2.9) on $`|\sigma _1`$ with $`n=m=0`$, we see that no term contributes on the left-hand side, so that using (2.15) and $`L_0|\sigma _1=h_1|\sigma _1`$ we obtain
$$\frac{2hh_1}{c}=\frac{1}{8},$$
$`(2.16)`$
which is a special case of (1.3) when $`c=c_{3,p}`$ (and recall that $`h=h_{2,1}`$). In other words, with $`c=c_{3,p}`$, this relation fixes the value of $`h_1`$ to $`h_{1,2}`$.
As previously pointed out, within the framework of this reformulation of the $`(3,p)`$ models in terms of the $`\varphi `$-algebra, a highest-weight $`\varphi `$-module is generically a combination of two Virasoro highest-weight modules (and this holds true irrespectively of the parity of $`p`$). Indeed, take for instance the vacuum state $`|0=|\sigma _0=|\varphi _{1,1}`$; its first descendant will be $`\varphi _h|0`$, which is itself the Virasoro highest-weight state $`|\varphi _{1,p1}=|\varphi _{2,1}`$. More generally, the states $`|\varphi _{1,s}`$ and $`|\varphi _{1,ps}`$ will be combined into a single module since $`|\varphi _{1,ps}`$ is a descendant of $`|\varphi _{1,s}`$ \- cf. (1.11). Note also that
$$𝒮|\varphi _{1,ps}=(1)^p|\varphi _{1,ps},$$
$`(2.17)`$
still with the understanding that $`s=\mathrm{}+1<p/2`$. The only case for which the $`\varphi `$-module reduces to a single Virasoro module is when $`p`$ is even and $`s=p/2`$.
2.2. The $`(3,5)`$ model
Let us first demonstrate, in a very explicit way, the necessity of the operator $`𝒮`$ anticommuting with $`\varphi `$. For this, we evaluate $`\varphi _{\frac{1}{4}}\varphi _{\frac{1}{4}}\varphi _{\frac{3}{4}}|0`$ in two different ways, symbolically written as:<sup>5</sup> This is an example of a mode-formulated associativity computation – cf. the appendix for more detail.
$$\underset{\mathrm{I}_{1,0,1}}{\underset{}{\varphi _{\frac{1}{4}}\varphi _{\frac{1}{4}}}}\varphi _{\frac{3}{4}}|0=\varphi _{\frac{1}{4}}\underset{\mathrm{I}_{1,0,0}}{\underset{}{\varphi _{\frac{1}{4}}\varphi _{\frac{3}{4}}}}|0,$$
$`(2.18)`$
using the notation introduced in (2.10). This relation means that we commute the first two terms using (2.6) with $`n=1,m=0,\mathrm{}=1`$ ($`\mathrm{}=1`$ because we act on a state with $`\mathrm{}=1`$, namely $`\varphi _{\frac{3}{4}}|0`$) and we compare this with the result of commuting the second and third terms using again (2.6) but now with $`n=1,m=0,\mathrm{}=0`$. That leads to
$$\varphi _{\frac{3}{4}}\varphi _{\frac{3}{4}}\varphi _{\frac{3}{4}}|0+𝒮\varphi _{\frac{3}{4}}|0=0,$$
$`(2.19)`$
(i.e., $`\varphi _{\frac{1}{4}}\varphi _{\frac{3}{4}}|0=0`$ follows directly form $`\mathrm{I}_{1,0,0}`$ and this is in agreement with the absence of a level-one descendant in the vacuum module: $`L_1|0=0`$). Next we use $`\mathrm{I}_{0,0,0}`$ to obtain
$$\varphi _{\frac{3}{4}}\varphi _{\frac{3}{4}}|0=𝒮|0.$$
$`(2.20)`$
The relation (2.19) becomes then
$$\left[\varphi _{\frac{3}{4}}𝒮+𝒮\varphi _{\frac{3}{4}}\right]|0=0,$$
$`(2.21)`$
which is precisely what we wanted to establish: $`𝒮`$ anticommutes with the modes of $`\varphi `$. On the other hand, there is no way of fixing the eigenvalue of $`𝒮`$ on $`|0`$ and it is thus chosen to be 1.
The central charge is obtained by evaluating
$$\underset{\mathrm{II}_{1,0,1}}{\underset{}{\varphi _{\frac{1}{4}}\varphi _{\frac{1}{4}}}}\varphi _{\frac{3}{4}}|0=\varphi _{\frac{1}{4}}\underset{\mathrm{II}_{1,0,0}}{\underset{}{\varphi _{\frac{1}{4}}\varphi _{\frac{3}{4}}}}|0.$$
$`(2.22)`$
This leads to
$$\frac{3}{8c}(5c+3)\varphi _{\frac{3}{4}}|0=0,$$
$`(2.23)`$
fixing $`c=3/5`$ as expected. Associativity is further tested by computing the central charge as
$$\underset{\mathrm{II}_{0,1,1}}{\underset{}{\varphi _{\frac{3}{4}}\varphi _{\frac{3}{4}}}}\varphi _{\frac{3}{4}}|0=\varphi _{\frac{3}{4}}\underset{\mathrm{II}_{0,0,0}}{\underset{}{\varphi _{\frac{3}{4}}\varphi _{\frac{3}{4}}}}|0.$$
$`(2.24)`$
This gives the equivalent result:
$$\frac{1}{8c}(7c+9)\varphi _{\frac{3}{4}}|0=\varphi _{\frac{3}{4}}|0.$$
$`(2.25)`$
Consider now the spectrum of the model. Since $`p=5`$, the bound (2.14) yields $`0\mathrm{}1`$. But we have already obtained the general dimension of the primary field $`\sigma _1`$ in (2.16). With $`p=5`$ and $`c=3/5`$, this yields $`h_1=1/20`$, which identifies $`\sigma _1`$ to $`\varphi _{1,2}`$. Its first descendant is
$$\varphi _{\frac{3}{4}+\frac{1}{2}}|\sigma _1=\varphi _{\frac{1}{4}}|\sigma _1|\varphi _{1,3},$$
$`(2.26)`$
of dimension $`1/5`$. With $`\sigma _0\varphi _{1,1}`$ and
$$\varphi _{\frac{3}{4}}|\sigma _0|\varphi _{1,4},$$
$`(2.27)`$
with dimension $`3/4`$, the spectrum of the $`(3,5)`$ model is completely recovered.
Let us briefly comment on the origin of the operator $`𝒮`$ in the context of the graded parafermionic models, with coset representation $`\widehat{osp}(1,2)_k/\widehat{u}(1)`$ . The $`(3,5)`$ model corresponds to $`k=1`$ . Let $`\psi _{\frac{1}{2}}`$ be the fundamental parafermionic field of dimension $`11/4k`$, which satisfies $`(\psi _{\frac{1}{2}})^{2k}I`$. Denote by $`\psi _1`$ the parafermion of dimension $`11/k`$. We have $`(\psi _{\frac{1}{2}})^2\psi _1`$. For $`k=1`$ however, $`\psi _1I`$. But this is true up to a zero mode. Indeed, if we denote by $`B`$ and $`A`$ the respective modes of $`\psi _{\frac{1}{2}}`$ and $`\psi _1`$, then we have
$$B_{\frac{1}{4}}A_1|0=0=[A_0B_{\frac{3}{4}}+B_{\frac{3}{4}}A_0]|0.$$
$`(2.28)`$
We thus recover (2.21) with $`B_{\frac{3}{4}}\varphi _{\frac{3}{4}}`$ and $`A_0𝒮`$.
2.3. The $`(3,7)`$ model
In the present case, the central charge is readily computed from
$$\underset{\mathrm{II}_{0,1,1}}{\underset{}{\varphi _{\frac{1}{4}}\varphi _{\frac{1}{4}}}}\varphi _{\frac{5}{4}}|0=\varphi _{\frac{1}{4}}\underset{\mathrm{II}_{1,0,0}}{\underset{}{\varphi _{\frac{1}{4}}\varphi _{\frac{5}{4}}}}|0,$$
$`(2.29)`$
leading to
$$\frac{1}{16c}(7c+25)\varphi _{\frac{5}{4}}|0=0,$$
$`(2.30)`$
with solution $`c=25/7`$.
The most interesting aspect to consider here, compared to the previous case, is the determination of the spectrum. In the $`(3,5)`$ case, we had two primary field: $`\sigma _0=I`$ and $`\sigma _1`$, whose dimension are directly determined by the commutation relations. With $`p=7`$, we have three primary fields: $`\sigma _0,\sigma _1`$ and $`\sigma _2`$. Again the dimension of $`\sigma _1`$ results from (2.9): $`h_1=5/28(=h_{1,2})`$. Here the difficulty lies in the determination of $`h_2`$, the dimension of $`\sigma _2`$, which does not follow from a direct application of the generalized commutation relations on $`|\sigma _2`$ . This dimension can be fixed by associativity, however. For instance, by comparing :
$$\underset{\mathrm{II}_{1,2,3}}{\underset{}{\varphi _{\frac{1}{4}}\varphi _{\frac{1}{4}}}}\varphi _{\frac{1}{4}}|\sigma _2=\varphi _{\frac{1}{4}}\underset{\mathrm{II}_{1,1,2}}{\underset{}{\varphi _{\frac{1}{4}}\varphi _{\frac{1}{4}}}}|\sigma _2,$$
$`(2.31)`$
we get
$$\left\{\frac{1}{16}\frac{5}{4c}\left(h_2+\frac{1}{4}\right)\right\}\varphi _{\frac{1}{4}}|\sigma _2=\frac{5h_2}{2c}\varphi _{\frac{1}{4}}|\sigma _2,$$
$`(2.32)`$
with solution $`h_2=1/7`$. We have thus
| | $`|\sigma _0|\varphi _{1,1}`$ | $`|\sigma _1|\varphi _{1,2}`$ | $`|\sigma _2|\varphi _{1,3}`$ |
| --- | --- | --- | --- |
| | $`\varphi _{\frac{5}{4}}|\sigma _0|\varphi _{1,6}`$ | $`\varphi _{\frac{3}{4}}|\sigma _1|\varphi _{1,5}`$ | $`\varphi _{\frac{1}{4}}|\sigma _2|\varphi _{1,4},`$ |
$`(2.33)`$
and this complete the analysis of the $`(3,7)`$ model.
3. The $`(3,p)`$ quasi-particle basis
3.1. The general form of the basis
Let us now turn to a description of the basis of states for the $`(3,p)`$ models, with $`p`$ of both parities.
In the highest-weight module of $`|\sigma _{\mathrm{}}`$, the different states in the $`N`$-particle sector are of the form:
$$\varphi _{h+\frac{\mathrm{}}{2}+\frac{(N1)}{2}n_1}\mathrm{}\varphi _{h+\frac{\mathrm{}}{2}+\frac{(Ni)}{2}n_i}\mathrm{}\varphi _{h+\frac{\mathrm{}}{2}+\frac{1}{2}n_{N1}}\varphi _{h+\frac{\mathrm{}}{2}n_N}|\sigma _{\mathrm{}},$$
$`(3.1)`$
with some constraints on the $`n_i`$’s to be specified below. Note the cumulative contribution of the $`\varphi `$ charge. The highest-weight module is obtained by summing over all possible particle-sector $`N`$. The indices $`n_i`$ are integers for $`p`$ odd and alternate between integer and half-integer values when $`p`$ is even:
$$n_{N2i}\mathrm{},n_{N2i1}\mathrm{}+\frac{p1}{2}.$$
$`(3.2)`$
It is convenient to represent the string (3.1) by a sequence whose $`N`$ entries (in the $`N`$-particle sector) are minus the modes $`n_i`$’s, i.e.,
$$\varphi _{h+\frac{\mathrm{}}{2}+\frac{(N1)}{2}n_1}\mathrm{}\varphi _{h+\frac{\mathrm{}}{2}+\frac{1}{2}n_{N1}}\varphi _{h+\frac{\mathrm{}}{2}n_N}(n_1,\mathrm{},n_{N1},n_N),$$
$`(3.3)`$
We will define our quasi-particle basis in terms of a filling process on a ground state. As argued below, this ground state is described by the sequence
$$(\mathrm{}6,r+5,\mathrm{\hspace{0.17em}4},r+3,\mathrm{\hspace{0.17em}2},r+1,\mathrm{\hspace{0.17em}0}).$$
$`(3.4)`$
where we have introduced the notation
$$r=\frac{p5}{2}.$$
$`(3.5)`$
Notice the increase of 2 units at distance 2 (from left to right). The filling process amounts to add states corresponding to ordinary partitions on this ground state and sum over all particle sectors. The resulting sequences are not genuine partitions but merely ‘jagged partitions’ which satisfy
$$n_in_{i+1}r+1,n_in_{i+2}+2,n_N0.$$
$`(3.6)`$
These conditions are directly read off (3.4). Note that for $`r=0`$ ($`p=5`$), these are standard partitions with distinct parts. For $`r=1`$ ($`p=7`$), (3.6) describes standard partitions with a difference 2 condition between parts separated by the distance 2. In the generic case $`r>1`$, there is a possible increase of $`r1`$ between the part $`n_{N2i1}`$ and its right nearest neighbour $`n_{N2i}`$. Note that the difference 2 at distance 2 implies a further possible increase of $`r+1`$ between $`n_{N2i}`$ and $`n_{N2i+1}`$ (again this is a direct consequence of (3.4): for $`i=1`$, $`r+1`$ is the difference between $`2`$ and $`r+1`$). For instance, with $`p=14`$, so that $`r=9/2`$, our candidate-basis in the vacuum module is built on the ground state $`(\mathrm{}6,\mathrm{\hspace{0.17em}1}/2,\mathrm{\hspace{0.17em}4},3/2,\mathrm{\hspace{0.17em}2},7/2,\mathrm{\hspace{0.17em}0})`$, that is:
$$\mathrm{}\varphi _{3+3\mathrm{𝟔}}\varphi _{3+\frac{5}{2}\frac{\mathrm{𝟏}}{\mathrm{𝟐}}}\varphi _{3+2\mathrm{𝟒}}\varphi _{3+\frac{3}{2}+\frac{\mathrm{𝟑}}{\mathrm{𝟐}}}\varphi _{3+1\mathrm{𝟐}}\varphi _{3+\frac{1}{2}+\frac{\mathrm{𝟕}}{\mathrm{𝟐}}}\varphi _{3\mathrm{𝟎}}|0=\mathrm{}\varphi _6\varphi _1\varphi _5\varphi _0\varphi _4\varphi _1\varphi _3|0.$$
$`(3.7)`$
(Note in particular that $`\varphi _1\varphi _3|0L_2|0`$).<sup>6</sup> Recall that $`(3,14)W_3(3,7)`$ model, so that $`\varphi =W`$ in this context. Whether this basis can be lifted to a basis for the whole class of $`W_3(3,p)`$ models remains to be seen, however.
Arguments supporting (3.6) are presented in the following subsection.
The characterization of the basis of states is not quite complete since yet there is no way of distinguishing the different highest-weight modules. Indeed, for $`\mathrm{}2`$, a further restriction has to be imposed on the parts. The origin of this sort of boundary condition is simply that the lowest state in the 2-particle sector of the $`|\sigma _{\mathrm{}}`$ module must be of the form $`\varphi _{h+\frac{1}{2}+\frac{\mathrm{}}{2}m_0}\varphi _{h+\frac{\mathrm{}}{2}}|\sigma _{\mathrm{}}`$ of dimension $`2h\mathrm{}1/2+m_0`$. But this state has to be proportional to $`L_1|\sigma _{\mathrm{}}`$, which forces $`m_0=\mathrm{}r`$. Since $`m_0`$ is the lowest value that $`n_{N1}`$ can take, we have
$$n_{N1}\mathrm{}r.$$
$`(3.8)`$
Our hypothesis would be that this is the whole set of boundary conditions. This is conformed by state counting at low levels.
Let us illustrate the conditions (3.6) as well as the boundary condition (3.8) by listing the states of the $`\mathrm{}=4`$ module of the $`(3,11)`$ model at the first levels. We do this by writing the corresponding sequences $`(n_1,\mathrm{},n_N)`$. We also restrict ourself to the Virasoro module $`|\varphi _{1,5}`$ which means that we only consider states that contain an even number of $`\varphi `$ modes. The boundary condition (3.8) requires $`n_{N1}1`$. The descendant states up to level 6 are:
| | $`1:`$ | $`(1,0)`$ |
| --- | --- | --- |
| | $`2:`$ | $`(2,0),(1,1)`$ |
| | $`3:`$ | $`(3,0),(2,1)`$ |
| | $`4:`$ | $`(4,0),(3,1),(2,2),(1,3),(3,2,1,0)`$ |
| | $`5:`$ | $`(5,0),(4,1),(3,2),(2,3),(4,2,1,0),(3,3,1,0)`$ |
| | $`6:`$ | $`(6,0),(5,1),(4,2),(3,3),(2,4),(5,2,1,0),(4,3,1,0),(3,4,1,0),(4,2,2,0),(3,3,1,0).`$ |
$`(3.9)`$
For instance, the two sequences $`(2,4)`$ and $`(5,2,1,0)`$ correspond to the states
$$(2,4):\varphi _{\frac{1}{4}+\frac{1}{2}\mathrm{𝟐}}\varphi _{\frac{1}{4}\mathrm{𝟒}}|\sigma _4,(5,2,1,0):\varphi _{\frac{1}{4}+\frac{3}{2}\mathrm{𝟓}}\varphi _{\frac{1}{4}+1\mathrm{𝟐}}\varphi _{\frac{1}{4}+\frac{1}{2}\mathrm{𝟏}}\varphi _{\frac{1}{4}\mathrm{𝟎}}|\sigma _4,$$
$`(3.10)`$
which indeed both have level 6. The first state containing 6 modes $`\varphi `$ occurs at level 9 and it is associated to the sequence $`(5,4,3,2,1,0)`$.
3.2. The rationale for the condition (3.6)
The aim of this section is to justify the basis (3.6) from conformal field theory. Our argument is, to a large extend, maintained at a sketchy level but we expect that the main points can be recovered by a more rigorous analysis.
To investigate the structure of the $`\varphi `$-type quasi-particle basis, we consider the counting of independent states by treating successively the different particle sectors (recall that the ‘particle sector’ is the number of $`\varphi `$-modes acting on the highest-weight state). In the first stage of the analysis, we only invoke the generic features of the $`\varphi `$-algebra. Only in the later steps do we require this to be also a Virasoro minimal model.
We will stick to the vacuum module ($`\mathrm{}=0)`$ for simplicity. It will prove convenient to rewrite the states under the form
$$\mathrm{}\varphi _{hb_{N3}}\varphi _{ha_{N2}}\varphi _{hb_{N1}}\varphi _{ha_N}|0.$$
$`(3.11)`$
Observe the specific choice made for the sign in front of $`h`$ within the modes which is designed to facilitate the use of the commutation relations. In this notation, the cumulative charge of the $`\varphi `$-mode is absorbed into the $`a_i`$ and $`b_i`$ labels.
In the one-particle sector, the states are of the form $`\varphi _{hn}|0`$. The only constraint is the highest-weight condition (2.11) which forces $`n0`$.
Consider next the two-particle sector. The basic constraint here comes from the commutation relation (2.6). This relation certainly continues to hold true also for the generic version of the algebra under consideration, which does not affect the leading term of the OPE. <sup>7</sup> We stress that this commutation relation holds for both parities of $`p`$ if we assume that $`𝒮=I`$ for $`p`$ even (in which case the infinite sum truncates to a finite one). More generally, it holds for all value of $`p^{}`$ with the understanding that $`\varphi =\varphi _{p^{}1,1}`$ and $`h=h_{p^{}1,1}`$. It is rather immediate to see that all the states $`\varphi _{hm}\varphi _{hn}|0`$ which do not satisfy $`mn+2`$ can be reexpressed in terms of those that do satisfy this constraint.
Note that $`\varphi _{hm}\varphi _{hn}|0`$ with $`mn+2`$ is equivalent to $`\varphi _{h+1/2n_1}\varphi _{hn_2}|0`$ with $`n_1n_2r+1`$. This analysis of the two-particle sector can be directly transposed to a sequence of two adjacent modes within the bulk of a string of $`\varphi `$-modes, where the condition takes the form $`n_in_{i+1}r+1`$.
So far, we have succeeded in explaining the jagged nature of the sequences of $`n_i`$’s associated to the string of modes in descendant states. Next, we have to consider the independent states in the three-particle sector. We consider states of the form
$$\varphi _{hn^{}}\varphi _{hm}\varphi _{hn}|0,$$
$`(3.12)`$
and look for a constraint relating $`n^{}`$ to $`n`$. At this point, we move away from the study of the most general $`\varphi `$-type free basis and take into account the constraints coming from the fact that the models we consider are also Virasoro minimal models. As a result, every regular field appearing in the OPE $`\varphi (z)\varphi (w)`$ has to be rexpressible in terms of the energy-momentum tensor only.<sup>8</sup> To be plain, this immersion amounts to a reduction of the allowed fields appearing in the OPE $`\varphi (z)\varphi (w)`$. For instance, $`(GG)(w)`$ is the first field so removed for those superconformal models that are also Virasoro minimal models, namely for the $`(3,8)`$ and $`(4,5)`$ models. This means that when we compute the commutation relations between the $`\varphi `$ modes, we can pick up regular terms of arbitrary order and still only obtain Virasoro modes on the right-hand side of the commutation relations. This in turns has the obvious consequence that any state made of Virasoro modes can be expressed in terms of $`\varphi `$ modes (in even number) and vice-versa. This is the advantage we take into account in our next step.
We stress that our present objective is to determine some ordering condition on trilinear states by a recursive process. Given a target ordering condition, we need to show that states which do not satisfy this ordering condition can be expressed in terms of those which do satisfy it. With this in mind, we now introduce a simplifying trick.
Let us return for a moment to the left-hand side of the commutation relations (2.6) and (2.9). Given the way these relations are derived, we see that the more terms we pick up in the OPE $`\varphi (z)\varphi (w)`$, the higher is the gap between the left-most modes of the first sum and the right-most modes of the second sum. Knowing that we can pick up any regular term (since they are all expressible as combinations of $`T`$), let us suppose that we select a regular term of sufficiently high order such that the second sum contains only ordered states already considered so far in our analysis (that is, states that are more ordered than those under consideration at a given recursive step). For the purpose of the present discussion, we can thus ignore the contribution of this second sum. That leads to a simplified version of the commutation relations (where the equality is to be understood as modulo terms previously considered):
$$\underset{i=0}{\overset{\mathrm{}}{}}c^{(i)}\varphi _{hn^{}i}\varphi _{hm+i}=[L]_{n^{}m}+[LL]_{n^{}m}+\mathrm{}$$
$`(3.13)`$
where in this notation, the $`c^{(i)}`$ are unspecified constants and $`[L\mathrm{}L]_N`$ (with $`n`$ factors of $`L`$) represents a given linear combination of $`n`$ $`L`$ modes at level $`N`$. All the coefficients are fixed by the proper choice of commutation relations and the values of $`n^{}`$ and $`m`$.
Let us now assume at this point that the terms $`n^{}<n1`$ in (3.12) can be reorganized in terms of those with $`n^{}<n`$. Using (3.13) in (3.12) for $`n^{}=n1`$, we obtain:
$$\begin{array}{cc}& \varphi _{hn+1}\varphi _{hm}\varphi _{hn}|0\hfill \\ & =\left[\left([L]_{nm+1}+[LL]_{nm+1}+\mathrm{}\right)+\underset{i=1}{\overset{\mathrm{}}{}}c^{(i)}\varphi _{hni+1}\varphi _{hm+i}\right]\varphi _{hn}|0\hfill \\ & =a\varphi _{h2nm+1}|0+\varphi _{hn}\left([L]_{nm+1}+[LL]_{nm+1}+\mathrm{}\right)|0+\underset{i=1}{\overset{\mathrm{}}{}}c^{(i)}\varphi _{hni+1}\varphi _{hm+i}\varphi _{hn}|0\hfill \\ & =a\varphi _{h2nm+1}|0+\varphi _{hn}\underset{i=0}{\overset{\mathrm{}}{}}c^{(i)}\varphi _{hnm+i+1}\varphi _{hi}|0+\underset{i=1}{\overset{\mathrm{}}{}}c^{(i)}\varphi _{hni+1}\varphi _{hm+i}\varphi _{hn}|0\hfill \end{array}$$
$`(3.14)`$
where $`a`$ is some constant. In the second equality, $`\varphi _{hn}`$ has been commuted with the Virasoro modes, while in last one, the Virasoro modes have been reexpressed in terms of the $`\varphi `$ modes.
We see that most of the $`\varphi `$-trilinear states can be now written in the more ordered form $`n^{}n`$. If we get rid of the states already considered so far, we see that our initial state can be expressed in terms of the remaining non-ordered terms as:
$$\varphi _{hn+1}\varphi _{hm}\varphi _{hn}|0=[\varphi _{hn}\varphi _{hm+2}\varphi _{hn1}+\varphi _{hn1}\varphi _{hm+4}\varphi _{hn2}+\mathrm{}]|0.$$
$`(3.15)`$
The sum has to stop at some point as we can use the reordering of the bilinear terms in $`\varphi `$. Now we can repeat this for these new non-ordered states until we ultimately run out of such non-ordered terms. In other words, by starting with a given $`n^{}=ni`$, a recursive process allows us to obtain states of the form $`n^{}>ni`$ for $`i>0`$.
If we try to apply the same trick to eliminate states of the form $`n^{}>n`$, we end up with relations linking the $`n^{}>n`$ states to $`n^{}<n`$ states. By consistency, we expect these relations to be the same ones we have first obtained for $`n^{}<n`$.
Finally, if we consider the states with $`n^{}=n`$, we now notice that these states with $`n^{}=n`$ reappear along the derivation. Without explicitly calculating every coefficients in front of the states, we do not know whether some of the $`n^{}=n`$ states can be eliminated or not.
We can thus naturally expect that the conditions on the trilinear terms in $`\varphi `$ will lie somewhere in between the conditions $`n^{}n`$ and $`n^{}>n`$. Note that the above analysis can be applied to all $`(p,p^{})`$ models if we replace $`\varphi `$ by $`\varphi _{p^{}1,1}`$. From low-level state-counting checks, we can verify that the previous property is indeed verified for any $`(p^{},p)`$ model with $`p^{}3`$ (which ensures that $`\varphi _{p^{}1,1}\varphi _{1,1}`$). The most restricted case corresponds to $`p^{}=3`$ because the $`\varphi `$ singular vector arises at the lowest possible level, which is 2. This case lies exactly on the upper-bound constraint, that is, $`n^{}>n`$. In order to get the least restricted case, we have to examine the models for which the singular vector appears as deeply as possible in the module. As $`\varphi _{p^{}1,1}`$ has its first singular vector at level $`p^{}1`$, it corresponds to cases where $`p^{}`$ is large. It appears that the trilinear terms in those cases are actually more restricted than $`n^{}n`$, indicating that this lower-bound condition is not saturated.
We have just seen how the condition $`n^{}n`$ had to be respected. Now we will try to be slightly more explicit about the possible restrictions for $`n^{}=n`$. Let us consider (3.13) once again along with (3.12). Depending on the value of $`m`$, we do not always need to pick up the same regular term in order to eliminate the second sum in the commutation relations. The higher the value of $`m`$, the greater is the number of possibilities we have to write the commutation relations. Using all these different choices of commutation relations, we get a number of linearly independent relations taking all the form of the first equality in (3.14), with the coefficients before each terms differing from one choice to the other. If one could keep track of all the conditions coming form these different equalities, this could lead us to two possible outcomes. On the one hand, the result might be some intermediate ‘basis’ between the spanning set of states and the complete sought-for basis, upon which we would have to apply the restrictions coming from the removal of the first null-field in $`\varphi `$. On the other hand, the result could actually take care of all possible constraints. If this second possibility is the actual one, it would mean that the immersion of the $`\varphi `$-extended conformal field theories into the $`(p^{},p)`$ models fixes completely the Virasoro singular-vector structure.
We end up this section by displaying a sample computation supporting the later alternative. The example to be considered is the one for which $`\varphi `$ has dimension 3/2. The $`\varphi `$-algebra is thus a superconformal algebra. Considered also as a Virasoro minimal model (which embodies a truncation of the space of states), this is associative for two values of $`c`$, corresponding to the $`(3,8)`$ and $`(4,5)`$ models. We will show that in the former case, the state $`\varphi _{\frac{3}{2}}\varphi _{\frac{1}{2}}\varphi _{\frac{3}{2}}|0`$ can be eliminated.
The anticommutation relation for $`\varphi =G`$ takes the form
$$\{G_{n+\frac{1}{2}},G_{m\frac{1}{2}}\}=\frac{n(n+1)}{2}\delta _{n+m,0}+\frac{c}{3}L_{m+n}.$$
$`(3.16)`$
This is obtained by considering only the singular terms in the OPE $`G(z)G(w)`$. But if, instead, we pick up the contribution of more terms, in particular, up to and including the level-two descendants of $`T`$, we then get
$$\begin{array}{cc}& \underset{t0}{}C_2^{(t)}[G_{n+\frac{1}{2}t}G_{m\frac{1}{2}+t}+G_{m\frac{5}{2}t}G_{n+\frac{5}{2}+t}]=\frac{(n+3)!}{24(n1)!}\delta _{n+m,0}\hfill \\ & +\left[\frac{3(n+3)(m)}{2c}+\beta _1(n+m+3)(n+m+2)\right]L_{n+m}+\beta _2\underset{t0}{}[L_{2t}L_{n+m+2+t}+L_{n+m+1t}L_{1+t}],\hfill \end{array}$$
$`(3.17)`$
with $`\beta _1`$ and $`\beta _2`$ being the coefficients of $`T^{\prime \prime }`$ and $`(TT)`$ respectively. These constants are fixed by associativity to the values
$$\beta _1=\frac{9(c+1)}{4c(22+5c)}\mathrm{and}\beta _2=\frac{51}{2c(22+5c)}.$$
$`(3.18)`$
Further calculations show that the central charge must be restricted to the two values $`7/10`$ and $`21/4`$ as previously said. Using (3.16), we can write
$$G_{\frac{3}{2}}G_{\frac{1}{2}}G_{\frac{3}{2}}=\frac{3}{c}G_{\frac{3}{2}}L_2=\frac{3}{c}\left[G_{\frac{7}{2}}+L_2G_{\frac{3}{2}}\right].$$
$`(3.19)`$
Using now the commutation relations (3.17) with $`n=m=1`$ in order to get an expression for $`L_2`$ (as given by the first term on the right-hand side of (3.17)) acting on $`G_{3/2}`$, we find that
$$L_2G_{\frac{3}{2}}|0=\frac{3}{c}\left[\left(3\beta _2\frac{6}{c}\right)L_2G_{\frac{3}{2}}+\left(2\beta _2+5\frac{3}{2c}\right)G_{\frac{7}{2}}\right]|0.$$
$`(3.20)`$
For $`c=7/10`$, the second coefficient on the right-hand side vanishes while the first one reduces to 1; in other words, we end up with the identity $`L_2G_{\frac{3}{2}}|0=L_2G_{\frac{3}{2}}|0`$. Therefore, in that case, there is no relation between $`L_2G_{\frac{3}{2}}|0`$ and $`G_{\frac{7}{2}}|0`$. But for the other allowed value of $`c`$, which corresponds to that of the $`(3,8)`$ model, there is one such relation. (Actually, we have recovered here the expression for the $`\varphi _{2,1}`$ singular vector.) It implies that, in this precise case, we can eliminate the state $`G_{\frac{3}{2}}G_{\frac{1}{2}}G_{\frac{3}{2}}|0`$.
Higher order terms can be treated along these lines. But clearly, going deeper in the modules requires the computation of more and more terms in $`\varphi (z)\varphi (w)`$. These computations are thus rather complicated, in addition to be model-dependent. But they provide independent verifications of the stated conditions (3.6).
4. The $`(3,p)`$ fermionic-type characters
Our main assumption is that (3.6) provides a basis. This has been supported by heuristic considerations, some explicit computations and the comparison with known bases for small values of $`p`$. However, establishing (3.6) rigorously is a hard mathematical problem. We circumvent this by showing that these conditions do indeed lead us to the expected characters. More precisely, we demonstrate here that the enumeration of states subject to the conditions (3.6) together with the boundary condition (3.8), reproduces the known expressions for the $`(3,p)`$ characters in their fermionic form.
In view of enumerating all the states in a given module (i.e., constructing its character), it is convenient to transform the ground state into one for which the parts do not increase from left to right. Let us then add to the ground state (3.4) the staircase of $`(r1)`$-height step:
$$(\mathrm{},\mathrm{\hspace{0.17em}5}r4,\mathrm{\hspace{0.17em}4}r3,\mathrm{\hspace{0.17em}3}r2,\mathrm{\hspace{0.17em}2}r1,r,\mathrm{\hspace{0.17em}1}).$$
$`(4.1)`$
The shifted ground state is thus<sup>9</sup> We stress that this shifting is simply a relabeling of the ground state. Note also if the ground state (3.4) involves both integers and half-integers for $`p`$ even, the shifting process generates only integer parts since $`2r`$ is always integer.
$$(\mathrm{},4r+1,\mathrm{\hspace{0.17em}4}r+1,\mathrm{\hspace{0.17em}2}r+1,\mathrm{\hspace{0.17em}2}r+1,\mathrm{\hspace{0.17em}1},\mathrm{\hspace{0.17em}1}).$$
$`(4.2)`$
Partitions defined on this shifted ground state can be characterized as follows. These are partitions $`(\lambda _1,\lambda _2,\mathrm{},\lambda _N)`$ of length $`N`$,
$$\lambda _i\lambda _{i+1},\lambda _N1,$$
$`(4.3)`$
satisfying the supplementary condition
$$\lambda _i\lambda _{i+2}+2r.$$
$`(4.4)`$
These conditions follow from (3.6) and
$$\lambda _i=n_i+(Ni)(r1)+1.$$
$`(4.5)`$
In other words, parts in $`(\lambda _1,\mathrm{},\lambda _N)`$ that are separated by the distance 2 must then differ by at least a $`2r`$, with $`2r=p5`$.<sup>10</sup> If, instead, we subtract from the ground state (3.4) the staircase $`(\mathrm{},r+2,r+1,r)`$, it becomes $`(\mathrm{}0,r,0,r,0,r)`$. This is the ground state of jagged partitions of type $`0r`$ in the terminology of .. To these conditions, we need to add the boundary condition:
$$\lambda _{N1}\mathrm{}.$$
$`(4.6)`$
Let $`p_{r,\mathrm{}}(w,N)`$ be the number of partitions of length $`N`$ and weight $`w`$ (that is, $`w=_i\lambda _i`$) satisfying (4.4) and (4.6). Denote the corresponding generating function by
$$G_{r,\mathrm{}}(z,q)=\underset{w,N0}{}p_{r,\mathrm{}}(w,N)q^wz^N.$$
$`(4.7)`$
For $`0\mathrm{}k`$, this function can be obtained in closed form as a $`k`$-multiple sum:
$$G_{r,\mathrm{}}(z,q)=\underset{m_1,m_2,\mathrm{}m_k0}{}\frac{q^{𝐦\stackrel{~}{B}𝐦+\stackrel{~}{C}𝐦}z^{2(m_1+\mathrm{}+m_{k1})+m_k}}{(q)_{m_1}\mathrm{}(q)_{m_k}},$$
$`(4.8)`$
where $`k`$ is related to $`p`$ by
$$k=\left[\frac{p}{3}\right],$$
$`(4.9)`$
(where $`[x]`$ stands for the integer part of $`x`$) and it is understood that
$$𝐦\stackrel{~}{B}𝐦=\underset{i,j=1}{\overset{k}{}}m_i\stackrel{~}{B}_{ij}m_j,\stackrel{~}{C}𝐦=\underset{i=1}{\overset{k}{}}\stackrel{~}{C}_im_i.$$
$`(4.10)`$
The $`k\times k`$ symmetric matrix $`\stackrel{~}{B}`$ reads
$$\stackrel{~}{B}=\left(\begin{array}{ccccc}2r& 2r& \mathrm{}& 2r& r\\ 2r& 2r+1& \mathrm{}& 2r+1& r+\frac{1}{2}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 2r& 2r+1& \mathrm{}& 2r+k2& r1+\frac{k}{2}\\ r& r+\frac{1}{2}& \mathrm{}& r1+\frac{k}{2}& k1\end{array}\right),$$
$`(4.11)`$
while the vector $`\stackrel{~}{C}`$ takes the form
$$\stackrel{~}{C}_j=2r+j+1+\mathrm{max}(\mathrm{}j,0)\mathrm{for}j<k\mathrm{and}\stackrel{~}{C}_k=k+2.$$
$`(4.12)`$
Finally, in the denominator of (4.8), we made use of the notation
$$(q)_n=\underset{i=1}{\overset{n}{}}(1q^i).$$
$`(4.13)`$
The generating function (4.8) for $`\mathrm{}=0`$ (which turns out to hold also for $`\mathrm{}=1`$) has been found in (cf. Theorem 5.14 with $`M_i,N\mathrm{}`$). For $`\mathrm{}0`$ but $`\mathrm{}k`$, the boundary condition (4.6) induces a modification by terms linear in the $`m_i`$’s and these are easily fixed by looking at the lowest partitions in low ($`N=1,2`$) particle sectors. This is taken care by the second term in $`C_j`$, as demonstrated in appendix B. The latter is a verification proof and it breaks for $`\mathrm{}>k`$. A further analysis of these boundary terms is presented in appendix C. It is shown there that $`G_{r,\mathrm{}}`$ can be recovered from $`G_{r,0}`$ recursively, for all $`1\mathrm{}p/21`$. However, it appears that it is only for $`1\mathrm{}k`$ that the $`G_{r,\mathrm{}}`$ can be reconstructed in closed form, as a single fermionic multisum.
Our goal is to count partitions built on the ground state (3.4) subject to (3.8) and weight them by their proper conformal dimension, constructing thereby the character of the $`|\sigma _{\mathrm{}}`$ module. We thus start with $`G_{r,\mathrm{}}(z,q)`$ and enforce $`z^N`$ to be equal to a certain power of $`q`$ adjusted in order to: 1- undo the ground state shifting by taking out the staircase contribution (4.1), whose weight is denoted $`w_{\mathrm{stair}}`$; and 2- add the fractional part, with weight $`w_{\mathrm{frac}}`$, to correct for the fact that (3.4) does not take into account the fractional part of the $`\varphi `$-modes. These numbers are easily computed. On the one hand, the staircase $`((N1)r(N2),\mathrm{},2r1,r,1)`$ has weight
$$w_{\mathrm{stair}}=\frac{N}{2}\left[(N1)(r1)+2\right].$$
$`(4.14)`$
On the other hand, the fractional part has the following dimension
$$w_{\mathrm{frac}}=N\left[h\frac{\mathrm{}}{2}\frac{(N1)}{4}\right]=\frac{N}{4}\left(2r+42\mathrm{}N\right),$$
$`(4.15)`$
where the third term in the square bracket comes from the cumulative contribution of the $`\varphi `$ charges. We thus replace $`z^N`$ in (4.8) by $`q^{w_{\mathrm{frac}}w_{\mathrm{stair}}}`$, where
$$N=2(m_1+\mathrm{}+m_{k1})+m_k.$$
$`(4.16)`$
This leads to the following expression for the character $`\widehat{\chi }_{\mathrm{}}`$ (normalized such that its leading $`q`$ power is $`q^0`$, hence the hat),
$$\widehat{\chi }_{\mathrm{}}(q)=\underset{m_1,m_2,\mathrm{}m_k0}{}\frac{q^{𝐦B𝐦+C𝐦}}{(q)_{m_1}\mathrm{}(q)_{m_k}},$$
$`(4.17)`$
where, with $`1i,jk1`$
$$B_{ij}=\mathrm{min}(i,j),B_{jk}=B_{kj}=j/2,B_{kk}=\frac{k+1ϵ}{4},$$
$`(4.18)`$
and $`C`$ reads
$$C_j=\mathrm{max}(j\mathrm{},0),C_k=\frac{k1+ϵ\mathrm{}}{2},$$
$`(4.19)`$
with $`ϵ=0,1`$ defined by
$$p=3k+1+ϵ.$$
$`(4.20)`$
In terms of the Virasoro characters, $`\widehat{\chi }_{\mathrm{}}`$ decomposes as follows
$$\widehat{\chi }_{s1}(q)=q^{h_{1,s}+c/24}\left[\chi _{1,s}^{\mathrm{Vir}}(q)+q^{h_{1,ps}h_{1,s}}\chi _{1,ps}^{\mathrm{Vir}}(q)\right].$$
$`(4.21)`$
Recall that in our construction, the Virasoro primary field $`\varphi _{1,ps}`$ (for $`p>2s`$) is a $`\varphi `$-descendant of $`\varphi _{1,s}`$. The two Virasoro characters can be separated by the parity of $`N`$, which is the same as that of $`m_k`$: with $`m_k`$ even (odd) , we obtain $`\chi _{1,s}`$ ($`\chi _{1,ps}`$ respectively).
We recover thus the fermionic sums given in in a form similar to that displayed here. Their complete proof is presented in and their reexpression in the above form is worked out in .
5. Conclusion
In this work, we have considered the reformulation of the minimal models $`(3,p)`$ in terms of the algebra spanned by $`\varphi \varphi _{2,1}`$ and defined by the OPE (1.6). The structure of this algebra differs somewhat according to the parity of $`p`$: for $`p`$ even, $`2h\mathrm{}_+`$ while for $`p`$ odd, $`h\mathrm{}_+\pm 1/4`$. In the latter case, the associativity analysis forces the introduction of a Witten-type operator anticommuting with $`\varphi `$. A similar operator (but presented differently) has been found in the spinon formulation of the $`\widehat{su}(2)_1`$ model . It seems to characterize non-local algebras with generators of dimension $`h\mathrm{}_+\pm 1/4`$
The Hilbert spaces (highest-weight states and their descendants) have been completely described in terms of the $`\varphi `$-algebra. In particular, the modules are described by the successive action of the lowering $`\varphi `$-modes subject to specific constraints. In the $`N`$-particle sector, with strings of lowering modes written in the form (1.13), these constraints are given in (1.14). The highest-weight states themselves are distinguished by the integer $`\mathrm{}`$ whose range is $`0\mathrm{}p/21`$. Moreover, the $`\mathrm{}`$-dependence of the descendant states is fully captured by the condition $`n_{N1}\mathrm{}r`$. The obtained basis agrees with those previously found for $`p=4,\mathrm{\hspace{0.17em}5}`$ and $`8`$ and the one derived in for $`\mathrm{}=0`$.
In absence of a complete argumentation underlying the derivation of this basis, our considerations have been supplemented by general arguments and explicit computations relying on the observation that the fine structure of the defining $`(3,p)`$ OPE $`\varphi (z)\varphi (w)`$ encodes the complete information on the models, including its quasi-particle basis. Note that this analysis does not mimic that of the $`(2,p)`$ and $`𝒮(2,4\kappa )`$ cases. In these cases, the spanning basis of states is first obtained and a set of restrictions, arising from the identity null field, is then imposed. The spanning set of states for the $`\varphi `$-algebra, on which one could impose constraints such as the level-two $`\varphi `$ null field, has not been found yet.
The simplest way of verifying the correctness of the basis of states is to derive the character of the irreducible module of $`|\sigma _{\mathrm{}}`$. This is obtained by enumerating all these states (1.13)-(1.14) and summing over all values of $`N`$. The character is explicitly given by (4.17) when $`0\mathrm{}[p/3]`$. This agrees with the known fermionic form of these characters . Although for $`[p/3]<\mathrm{}p/21`$, the character has not been found in closed form, it is shown in appendix C how it can be obtained recursively form $`G_{r,0}`$. But we stress that the validity of the boundary term in this range has been tested by writting explicitly the states at the first few levels of various modules and comparing their enumeration with that given by the usual bosonic formula.
The $`(3,p)`$ models, due to the presence of an extra symmetry generator, are somewhat similar to the superconformal models. For the special class of $`𝒮(2,4\kappa )`$ models, we can either choose to write the quasi-particle basis either in terms of the modes of $`G`$ together with the Virasoro modes or solely in terms of the $`G`$ modes . This suggests that one could look for an alternative quasi-particle basis for the $`(3,p)`$ models, this one formulated in term of an ordered set of Virasoro modes acting on an ordered set of $`\varphi `$ modes as
$$L_{n_1}\mathrm{}L_{n_k}\varphi _{m_1}\mathrm{}\varphi _{m_k^{}}|0$$
$`(5.1)`$
with some constraints on the numbers $`n_i`$ and $`m_i`$. A basis of that type has indeed been found; this result will be presented elsewhere .
Another natural axis for extension follows from the observation that an analysis similar to the present one should be applicable to all extended algebras having the essential simplifying property of being single-channel.
Let us conclude by emphasizing the fact that there is a relatively small number of conformal field theories for which the fermionic characters are described in terms of a basis derived by intrinsic conformal field theoretical methods. We have already mentioned that this is so for the $`(2,p)`$ minimal models together with their superconformal analogues, the $`𝒮(2,4\kappa )`$ models . But there are few other examples, like the $`\widehat{su}(2)_k`$ models , particular higher-rank WZW models , the parafermionic models , their graded version and some higher-rank formulations . The present analysis is a step toward the addition of a further example to this list, the $`(3,p)`$ minimal models.
Appendix A. Associativity and Jacobi identities
The associativity conditions for the symmetry generators of an extended conformal algebra are sometimes loosely viewed as being equivalent to the Jacobi identities for the mode-generators. For fields with (half-)integer dimension that associativity implies the Jacobi identity can indeed be derived in a rather direct way.<sup>11</sup> Given three integer-dimension operators $`A,B,C`$ with ordinary commutation relations, let us consider the action of $`A_nB_mC_p`$ on an arbitrary state $`|h`$. The associativity requirement retranscribed at the level of modes implies that this state can be evaluated by commuting the first two terms or the last two ones without affecting the result. This in turn implies the usual form of the Jacobi identity as we now show. Commuting $`A`$ and $`B`$ and reexpressing the result in terms of the state $`C_pB_mA_n|h`$ yields
$$\left([[A_n,B_m],C_p]+C_p[A_n,B_m]+B_m[A_n,C_p]+[B_m,C_p]A_n+C_pB_mA_n\right)|h.$$
Commuting $`B`$ and $`C`$ and singling out again the state $`C_pB_mA_n|h`$ gives
$$\left([A_n,[B_m,C_p]]+[B_m,C_p]A_n+[A_n,C_p]B_m+C_p[A_n,B_m]+C_pB_mA_n\right)|h.$$
The comparison of these two expressions implies the identity
$$\left([[A_n,B_m],C_p]+[[C_p,A_n],B_m]+[[B_n,C_p],A_n]\right)|h=0.$$
This is the Jacobi identity; more precisely, this is the Jacobi identity modulo a singular vector of $`|h`$. For fields with half-integer dimension, this is also true but with an appropriate graded version of the Jacobi identity. But the reverse is not true. As the following considerations will illustrate, associativity contains more information than the mere Jacobi identities for the modes.
Consider a free fermion, whose OPE and mode decomposition (in the NS sector) read:
$$\psi (z)\psi (w)\frac{1}{zw},\psi (z)=\underset{n\mathrm{}+1/2}{}b_nz^{n1/2}.$$
$`(\text{A.}1)`$
By evaluating the integral
$$\frac{1}{(2\pi i)^2}_0𝑑w_w𝑑zz^{n1/2}w^{m1/2}\psi (z)\psi (w),$$
$`(\text{A.}2)`$
we find the usual anti-commutation relations:
$$\{b_n,b_m\}=\delta _{m+n,0}.$$
$`(\text{A.}3)`$
Considering this anticommutator together with the Virasoro commutation relations and
$$[L_n,b_m]=\left(\frac{n}{2}+m\right)b_{n+m},$$
$`(\text{A.}4)`$
it is simple to convince oneself that the central charge is not fixed by the Jacobi identity. However, if we consider the following version of the OPE
$$\psi (z)\psi (w)=\frac{1}{zw}+\frac{zw}{c}T(w)+\mathrm{},$$
$`(\text{A.}5)`$
(given that the $`\beta ^{(2)}`$ coefficient is $`2h_\psi /c=1/c`$) and the integral
$$\frac{1}{(2\pi i)^2}_0𝑑w_w𝑑z\frac{z^{n+1/2}w^{m+1/2}}{(zw)^2}\psi (z)\psi (w),$$
$`(\text{A.}6)`$
we get the following generalized relations:
$$\underset{l0}{}l[b_{n1l}b_{m+1+l}+b_{m1l}b_{n+1+l}]=\frac{(n1/2)(n+1/2)}{2}\delta _{n+m,0}+\frac{1}{c}L_{n+m},$$
$`(\text{A.}7)`$
or equivalently (cf. , App. C)
$$\underset{l0}{}(l+1)[b_{n3/21}b_{m+3/2+l}+b_{m1/21}b_{n+1/2+l}]=\frac{n(n1)}{2}\delta _{n+m,0}+\frac{1}{c}L_{n+m}.$$
$`(\text{A.}8)`$
These relations capture the expression of the energy-momentum of the free fermion in terms of modes (but as a function of the yet-to-be-fixed central charge). In other words, the relation between $`T`$ and $`\psi `$ is already coded in the OPE and the conformal invariance.<sup>12</sup> This, of course, is a totally standard statement since the OPE (A.5) can also be written as
$$\psi (z)\psi (w)=\frac{1}{zw}+(\psi (z)\psi (w))=\frac{1}{zw}+(zw)(\psi \psi )(w)+\mathrm{},$$
from which we read that
$$T(w)=c(\psi \psi )(w)=c(\psi \psi )(w),$$
and this corresponds to the usual expression when $`c=1/2`$. Note that the four-point function
$$\psi _1(z_1)\psi _2(z_2)\psi _3(z_3)\psi _4(z_4)=\frac{1}{z_{12}z_{34}}\frac{1}{z_{13}z_{24}}+\frac{1}{z_{14}z_{23}},$$
$`(\text{A.}9)`$
which is evaluated solely from the knowledge of the singular terms (through meromorphicity, cf. ), contains the information on the central charge when viewed in the light of the OPE (A.5). It is extracted by evaluating the correlator in the limit $`z_1z_2`$. This does not contradict the previous conclusion because the correlation function encompasses the whole content of the involved OPE.
The free-boson theory offers another simple illustration of the gap between the information obtained from associativity and the Jacobi identities. The OPE
$$i\phi (z)i\phi (w)=\frac{1}{(zw)^2}+\frac{2}{c}T(w)+\mathrm{},$$
$`(\text{A.}10)`$
and mode decomposition
$$i\phi =\underset{n}{}a_nz^{n1},$$
$`(\text{A.}11)`$
lead to the standard mode commutation relation as
$$[a_n,a_m]=\frac{1}{(2\pi i)^2}_0𝑑w_w𝑑zz^nw^mi\phi (z)i\phi (w)=n\delta _{n+m,0},$$
$`(\text{A.}12)`$
for which the Jacobi identity is trivial. If we consider instead the integral
$$\frac{1}{(2\pi i)^2}_o𝑑w_w𝑑z\frac{z^nw^{m+1}}{(zw)}i\phi (z)i\phi (w),$$
$`(\text{A.}13)`$
we end up with a very different form of the commutation relations, namely,
$$\underset{l0}{}[a_{nl1}a_{m+1+l}+a_{ml}a_{n+l}]=\frac{n(n1)}{2}\delta _{n+m,0}+\frac{2}{c}L_{n+m}.$$
$`(\text{A.}14)`$
It is easily checked that $`c`$ is fixed to $`1`$ by associativity. This last expression is thus seen to be the mode expression of the energy-momentum of the free boson, i.e.,
$$L_n=\underset{l\mathrm{}}{}a_la_{nl},(n0)L_0=\underset{l0}{}a_la_l+\frac{1}{2}a_0^2.$$
$`(\text{A.}15)`$
This observation is of course true in general: picking up the energy-momentum as the single pole, yields directly the expression of the Virasoro modes in terms of the conserved-current ones.
These simple considerations show that the OPE contains more information than a particular form of commutation relation derived from it. To recover the complete information which is contained in the OPE, we need to consider the infinite family of commutation relations that follows from evaluating (A.6) with $`(zw)^2(zw)^p`$, for all values of $`p`$. In practice, however, only few values of $`p`$ should be sufficient.
The associativity of the four-point functions $`ABCD`$ is equivalent to the statement that the state $`A_{\mathrm{}}B_mC_n|h`$ where $`h`$ is the dimension of $`D`$, is independent of the way it is evaluated. Generically, the resulting constraints are independent of the state $`|h`$ (unless an equality is true modulo a singular vector) and in practice it can be replaced by the vacuum. To formulate the mode version of the associativity constraints, we introduce the notation
$$\underset{p}{\underset{}{AB}}C[A,B]_pC+BAC,$$
$`(\text{A.}16)`$
where $`[A,B]_p`$ stands for the commutator evaluated in terms of the generalized commutation relations that follow from evaluating the integral:
$$\frac{1}{(2\pi i)^2}_0𝑑w_w𝑑z\frac{z^n^{}w^m^{}}{(zw)^p}A(z)B(w),$$
$`(\text{A.}17)`$
(the values of $`n^{}`$ and $`m^{}`$ being adapted to the choice of $`p`$ in order to recover a final commutator in standard form – cf. (A.2) vs (A.6) and (A.12) vs (A.13)). Now, mode-associativity boils down to the following conditions:
$$\underset{p}{\underset{}{AB}}C=\underset{q}{\underset{}{AB}}C=A\underset{p}{\underset{}{BC}}.$$
$`(\text{A.}18)`$
In principle, these conditions should be tested for all values of $`p`$ and $`q`$ and all combinations of modes. However, in practice, a small number of computations of this type are needed to fix the whole structure of the models under consideration.
As a simple illustration, let us show how we can fix the central charge of the free-fermion theory by enforcing the mode-associativity of $`b_{1/2}b_{1/2}b_{1/2}|0`$. For this we compare
$$b_{1/2}\underset{p=0}{\underset{}{b_{1/2}b_{1/2}}}|0=b_{1/2}|0,$$
$`(\text{A.}19)`$
(the $`p=0`$ commutator being (A.2)) to
$$\underset{p=2}{\underset{}{b_{1/2}b_{1/2}}}b_{1/2}|0=\frac{1}{c}L_0b_{1/2}|0=\frac{1}{2c}b_{1/2}|0,$$
$`(\text{A.}20)`$
(the $`p=2`$ commutator being (A.6)) to find that $`c=1/2`$. Note that in this case, we could also compute the central charge by comparing the last result with
$$b_{1/2}\underset{p=2}{\underset{}{b_{1/2}b_{1/2}}}|0=b_{1/2}|0,$$
$`(\text{A.}21)`$
(in the first case, we generate a term proportional to a Virasoro mode, hence containing $`c`$, and in the second case the remaining contribution is the delta term, independent of $`c`$).
Appendix B. Boundary terms in the generating function
Our goal is to derive the modification of the generating function $`G_{r,0}(z,q)`$ (4.8), that counts the partitions $`(\lambda _1,\mathrm{},\lambda _N)`$ with $`\lambda _i\lambda _{i+2}+2r`$ with $`\lambda _N1`$ , which results from the further condition $`\lambda _{N1}\mathrm{}`$. We start with the assumption that the boundary terms are represented by linear factors in the exponent, i.e., are accounted by a correction of the form $`q^{_{j=1}^km_ja_j}`$. Denote the modified form as $`G_{r,\mathrm{}}(z,q)`$. Since the boundary condition is well localized i.e., it concerns only the second term at the right, it suffices to consider the $`N=1,2`$ sectors to fix the $`a_j`$. Recall that $`N=2_{j=1}^{k1}m_j+m_k`$. Therefore, $`N=1`$ corresponds to $`m_k=1`$ and all the others $`m_j=0`$. The generating function $`G_{r,0}`$ multiplied by $`q^{a_k}`$ reads
$$G_{r,\mathrm{}}(q,1)=\frac{q^{1+a_k}}{(q)_1}(N=1).$$
$`(\text{B.}1)`$
But the partitions that are to be counted are simply those with a single part and their generating function is $`q/(q)_1`$. This fixes $`a_k=0`$.
Consider next $`N=2`$ which requires either that a single mode $`m_j=1`$ for $`1jk1`$ with all other modes zero or that $`m_k=2`$, again with all other modes vanishing. This results into
$$G_{r,\mathrm{}}(q,1)=\frac{1}{(q)_1}\underset{j=1}{\overset{k1}{}}q^{2j+a_j}+\frac{q^{2k}}{(q)_2}(N=2)$$
$`(\text{B.}2)`$
In order to fix the $`a_j`$, we must determine the generating function enumerating partitions with two parts $`(\lambda _1,\lambda _2)`$ and satisfying
$$\lambda _1\lambda _21\mathrm{and}\lambda _1\mathrm{}.$$
$`(\text{B.}3)`$
This can be done by means of the MacMahon method (cf. vol 2 Sect. VIII, chap. 1 and see also , Sect.. 11.2), by projecting the following expression
$$\frac{a_3^{\mathrm{}}a_2^1}{(1a_1a_3q)(1a_2q/a_1)}=\underset{\lambda _1,\lambda _20}{}a_1^{\lambda _1\lambda _2}a_2^{\lambda _21}a_3^{\lambda _1\mathrm{}}q^{\lambda _1+\lambda _2},$$
$`(\text{B.}4)`$
onto positive powers of the $`a_i`$’s, ensuring thereby the three inequalities: $`\lambda _1\lambda _2`$, $`\lambda _21`$ and $`\lambda _1\mathrm{}`$. It is convenient to introduce the MacMahon projection symbol $`\mathrm{\Omega }`$, defined by
$$\underset{}{\overset{a}{\Omega }}\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}c_na^n=\underset{n0}{}c_na^n|_{a=1}=\underset{n0}{}c_n$$
$`(\text{B.}5)`$
and make use of identities of the following type:
$$\underset{}{\overset{a}{\Omega }}\frac{1}{(1aq)(1a^1q)}=\underset{}{\overset{a}{\Omega }}\frac{1}{(1q^2)}\left(\frac{1}{1aq}+\frac{a^1q}{1a^1q}\right)=\frac{1}{(1q)(1q^2)}.$$
$`(\text{B.}6)`$
The first two projections are rather direct
$$\underset{}{\overset{a_3}{\Omega }}\underset{}{\overset{a_2}{\Omega }}\underset{}{\overset{a_1}{\Omega }}\frac{a_3^{\mathrm{}}a_2^1}{(1a_1a_3q)(1a_2q/a_1)}=\underset{}{\overset{a_3}{\Omega }}\underset{}{\overset{a_2}{\Omega }}\frac{a_3^{\mathrm{}}a_2^1}{(1a_3q)(1a_2a_3q^2)}=\underset{}{\overset{a_3}{\Omega }}\frac{a_3^1\mathrm{}q^2}{(1a_3q)(1a_3q^2)},$$
$`(\text{B.}7)`$
and for the third one, we have
$$\begin{array}{cc}\hfill \underset{}{\overset{a}{\Omega }}\frac{a^1\mathrm{}q^2}{(1aq)(1aq^2)}& =\frac{a^1\mathrm{}q^2}{(1aq)(1aq^2)}\left(1(1aq)(1aq^2)\underset{p=0}{\overset{\mathrm{}2}{}}\underset{j=0}{\overset{p}{}}a^pq^{p+j}\right)|_{a=1}\hfill \\ & =\frac{q^{\mathrm{}+1}+q^{\mathrm{}+2}q^{2\mathrm{}+1}}{(1q)(1q^2)}\hfill \end{array}.$$
$`(\text{B.}8)`$
The $`a_j`$ are then fixed by comparing (B.2) and (B.8), the solution of which being
$$a_j=\mathrm{max}(\mathrm{}j,0),$$
$`(\text{B.}9)`$
as announced (cf. (4.12)). This is valid for $`0\mathrm{}k`$.
Appendix C. The analysis of boundary terms via recurrence relations for generating functions
We reconsider the construction of the generating functions for partitions $`\lambda =(\lambda _1,\mathrm{},\lambda _N)`$ into $`N`$ parts satisfying $`\lambda _i\lambda _{i+1}`$ and $`\lambda _i\lambda _{i+2}+2r`$, together with the boundary condition $`\lambda _{N1}\mathrm{}`$. The set of such partitions can be described schematically as (see e.g., )
$$\mathrm{}(2r+\mathrm{})(2r+1)(\mathrm{})(1)^+,$$
$`(\text{C.}1)`$
indicating that we build up these restricted partitions on the above ground state and the $`+`$ sign indicates the position from which we start the building up process (from right to left) by addition of ordinary partitions. We then use this pictural representation to write down the recurrence relation between sets with different boundary conditions:
$$\begin{array}{cc}\hfill \mathrm{}(2r+\mathrm{})(2r+1)(\mathrm{})(1)^+\mathrm{}(2r+\mathrm{}+1)(2r+1)(\mathrm{}+1)(1)^+& \\ \hfill =\mathrm{}(2r+\mathrm{})(2r+1)^+(\mathrm{})(1)+\mathrm{}(2r+\mathrm{})(2r+2)^+(\mathrm{})(2)& \\ \hfill +\mathrm{}+\mathrm{}(2r+\mathrm{})(2r+\mathrm{})^+(\mathrm{})(\mathrm{})& \end{array}$$
$`(\text{C.}2)`$
The difference on the left-hand side generates the set of partitions for which the penultimate part is $`\mathrm{}`$ and this set is then broken, on the right-hand side, into sets with prescribed values of the last two entries. Denote by $`p_{\mathrm{}}(w,N)`$ the number of partitions of weight $`w`$ (where $`w=\lambda _i`$) with $`N`$ parts in the set (C.1). The above recurrence relation can be translated into the following condition
$$p_{r,\mathrm{}}(w,N)p_{r,\mathrm{}+1}(w,N)=\underset{s=0}{\overset{\mathrm{}1}{}}p_{r,\mathrm{}s}(w(2r+s)(N2)\mathrm{}s1,N2)$$
$`(\text{C.}3)`$
On the left-hand side, we have used the observation that the set with the tail $`(\mathrm{})(s+1)`$ (for $`0s\mathrm{}1`$) is in one-to-one correspondence with the set obtained by deleting the last two parts $`(\mathrm{})(s+1)`$ and subtracting $`2r+s`$ from each of the $`N2`$ remaining parts, whose cardinality is thereby given by $`p_{r,\mathrm{}s}(w(2r+s)(N2)\mathrm{}s1,N2)`$. In terms of the generating function (4.7), the above recurrence relation implies:
$$G_{\mathrm{}+1}(z,q)=G_{\mathrm{}}(z,q)z^2\underset{s=0}{\overset{\mathrm{}1}{}}q^{\mathrm{}+s+1}G_\mathrm{}s(zq^{2r+s},q).$$
$`(\text{C.}4)`$
To this recurrence relation, we add the boundary condition $`G_1=G_0`$, with $`G_0`$ given in . In this way, we can construct $`G_{\mathrm{}}`$ recursively out of the known expression for $`G_0`$. In the following, we will only need the specialized version at $`z=1`$:
$$G_{r,\mathrm{}+1}(1,q)=G_{r,\mathrm{}}(1,q)\underset{s=0}{\overset{\mathrm{}1}{}}q^{\mathrm{}+s+1}G_{r,\mathrm{}s}(q^{2r+s},q).$$
$`(\text{C.}5)`$
Let us now show that for $`0\mathrm{}k`$, this relation leads to the expression already presented in (4.8), (4.11) and (4.12). Take some $`\mathrm{}+1k`$ and start by considering the difference
$$G_{r,\mathrm{}}^{(0)}(1,q)G_{r,\mathrm{}}(1,q)q^{\mathrm{}+1}G_{r,\mathrm{}}(q^{2r},q)$$
$`(\text{C.}6)`$
We first break the multisum expression for $`G_{\mathrm{}}(1,q)`$ into two parts: one with $`m_1=0`$ and the other with $`m_1>0`$. In the second sum, we redefine $`m_1=m_1^{}+1`$. The resulting expression has the same numerator as $`q^{\mathrm{}+1}G_{r,\mathrm{}}(q^{2r},q)`$ so that their difference is easily computed. Recombining the result with the first multisum corresponding to $`m_1=0`$ yields an expression, denoted $`G_{r,\mathrm{}}^{(0)}(1,q)`$ above, that is identical to $`G_{r,\mathrm{}}(1,q)`$ except that the linear coefficient of $`m_1`$ has been increased by 1. Next let us consider the difference between this expression $`G_{r,\mathrm{}}^{(0)}(1,q)`$ and the second term ($`s=1`$) of the sum of the right-hand side of (C.5):
$$G_{r,\mathrm{}}^{(1)}(1,q)G_{r,\mathrm{}}^{(0)}(1,q))q^{\mathrm{}+2}G_{r,\mathrm{}}(q^{2r+1},q).$$
$`(\text{C.}7)`$
The same manipulations but with $`m_1`$ replaced by $`m_2`$ gives back $`G_{r,\mathrm{}}^{(0)}(1,q)`$ except that the linear coefficient of $`m_2`$ has been increased by 1. Proceeding in this way by successively taking into account the different terms of the sum, we end up with the prescription that going from $`G_{r,\mathrm{}}`$ to $`G_{r,\mathrm{}+1}`$, amount to increase by 1 all the linear coefficients of $`m_1,\mathrm{},m_\mathrm{}1`$ and $`m_{\mathrm{}}`$ by 1 without modifying the other mode coefficients. This demonstrates that the difference between the linear term without boundary condition, pertaining to $`\mathrm{}=0,1`$, and those with $`\lambda _{N1}\mathrm{}`$, with $`1\mathrm{}k`$, is precisely given by the term max $`(\mathrm{}j)`$ in (4.12).
By construction, this recombination of the right-hand side of (C.5) works for $`\mathrm{}+1k`$. This bound can be justified by the argument presented in appendix B concerning the independence of the coefficient of $`m_k`$ upon $`\mathrm{}`$: there are thus no more coefficient available to reshuffle.To see explicitly how this recombination process breaks down for $`\mathrm{}>k`$, it suffices to consider $`\mathrm{}=3`$ and $`p=8`$ (so that $`k=2`$). This suggests that for $`\mathrm{}>k`$, the generating function $`G_{r,\mathrm{}}(z,q)`$ might not be expressible in terms of a single multiple sum.
ACKNOWLEDGMENTS
The work of PJ is supported by EPSRC and partially by the EC network EUCLID (contract number HPRN-CT-2002-00325), while that of PM is supported by NSERC.
REFERENCES
relax1.P. Jacob and P. Mathieu, Nucl. Phys. B 620 (2002) 351.
relax2.J. Lepowsky and M. Primc, Contemporary Mathematics 46 AMS, Providence, 1985.
relax3.J. Lepowsky and R.L. Wilson, Proc. Nat. Acad. Sci. USA 78 (1981) 7254.
relax4.B.L. Feigin, T. Nakanishi and H. Ooguri, Int. J. Mod. Phys. A7 Suppl. 1A (1992) 217.
relax5.R. Kedem, T.R. Klassen, B. M. McCoy and E. Melzer, Phys. Lett. B304 (1993) 263 and Phys. Lett. B307 (1993) 68; R. Kedem, B. M. McCoy and E. Melzer, in Recent progress in statistical mechanics and quantum field theory, ed. by P. Bouwknegt et al, World Scientific (1995) 195.
relax6.E. Melzer, Int. J. Mod. Phys. A9 (1994) 1115; A. Berkovich, Nucl. Phys. B241 (1984) 333.
relax7.A. Berkovich and B.M. McCoy, Lett. Math. Phys. 37 (1996) 49; A. Berkovich, B.M. McCoy and A. Schilling, Comm. Math Phys. 191 (1998) 329.
relax8.S.O. Warnaar, J. Stat. Phys. 82 (1996) 657 and J. Stat. Phys. 84 (1997) 49.
relax9.G.E. Andrews, R.J. Baxter and P.J. Forrester, J. Stat. Phys. 35 (1984) 193.
relax10.T. Welsh, Fermionic expressions for minimal model Virasoro characters, to appear in Memoirs of the American Mathematical Society, math.CO/0212154.
relax11.B. Feigin and A. V. Stoyanovsky, Quasi-particles models for the representations of Lie algebras and geometry of flag manifold, hep-th/9308079.
relax12.B.L. Feigin and E. Frenkel, Adv. Sov. Math. 16 (1993) 139.
relax13.A. Capelli, C. Itzykson and J.-B. Zuber, Comm. Math. Phys 13 (1987) 1.
relax14.G. Moore and N. Seiberg, Nucl. Phys B313 (1989) 16.
relax15.P. di Francesco and P. Mathieu, Phys. Lett. B270 (1992) 79.
relax16.A.B. Zamolodchikov and V.A. Fateev, Sov. Phys. JETP 43 (1985) 215.
relax17.D. Bernard, V. Pasquier and D. Serban , Nucl. Phys. B428 (1994) 612.
relax18.P. Bouwknegt, A. W.W. Ludwig and K. Schoutens, Phys. Lett. B338 (1994) 448.
relax19.P. Bouwknegt, A. W.W. Ludwig and K. Schoutens, Phys. Lett. B359 (1995) 304
relax20.A.L. Zamolodchikov, Theo. Math. Phys. 63 (1985) 1205.
relax21.A.A. Belavin, A.M. Polyakov and A.B. Zamolodchikov, Nucl. Phys. B241 (1984) 333.
relax22.Vl.S. Dotsenko and V.A. Fateev, Nucl. Phys. B235 (1984) 312; Nucl. Phys. B251 (1985) 691; Phys. Lett. 154B (1985) 291.
relax23.V.G. Knizhnik and A.B. Zamolodchikov, Nucl. Phys. B247 (1984) 83.
relax24.P. Bouwknegt, Phys. Lett. B207 (1988) 295.
relax25.E. Witten, Nucl. Phys B202 (1992) 253.
relax26.P. Jacob and P. Mathieu, Nucl. Phys. B630 (2002) 433
relax27.J. M. Camino, A. V. Ramallo and J. M. Sanchez de Santos, Nucl.Phys. B530 (1998) 715.
relax28.J.-F. Fortin, P. Jacob and P. Mathieu, J. Phys. A: Mat. Gen. 38 (2005) 1699.
relax29.A.G. Bytsko, J. Phys. A: Math. Gen. 32 (1999) 8045; A.G. Bytsko and A. Fring, Comm. Math. Phys. 209 (2000) 179.
relax30.G. Andrews, Pacific J. Math. 114 (1984) 267.
relax31.O. Foda and Y.-H.Quano, Int. J. Mod. Phys. A12 (1997) 1651.
relax32.B. Feigin, O. Foda and T. Welsh, Andrews-Gordon identities from combinations of Virasoro characters, math-ph/0504014.
relax33.B. Feigin, M. Jimbo and T. Miwa, Vertex operator algebra arising from the minimal series $`M(3,p)`$ and monomial basis, in MathPhys odyssey, Prog. Math. Phys., 23, Birkhäuser (2002) 179.
relax34.B. Feigin, M. Jimbo, T. Miwa, E. Mukhin and Y. Takayama, Set of rigged paths with Virasoro characters, math.QA/0506150.
relax35.L. Bégin, J.-F. Fortin, P. Jacob and P. Mathieu, Nucl. Phys B659 (2003) 365.
relax36.J.-F. Fortin, P. Jacob and P. Mathieu, Jagged partitions, Ramanujan J. , to appear, math.CO/0310079.
relax37.P. Jacob and P. Mathieu, A new quasi-particle basis for the $`M(3,p)`$ models, in preparation.
relax38. P. Bouwknegt. and K. Schoutens, Nucl. Phys. B 547 (1999) 501.
relax39.G. Georgiev, J. Pure Appl. Algebra 112 (1996) 247.
relax40.G. Georgiev, Combinatorial constructions of modules for infinite-dimensional Lie algebras, II. Parafermionic space, q-alg/9504024.
relax41.P. Jacob and P. Mathieu, Nucl. Phys. B 587 (2000) 514.
relax42.P. Jacob and P. Mathieu, J. Phys. A: Math. Gen. 34 (2001) 10141.
relax43.P. MacMahon, Combinatory analysis, 2 vols (1917,1918), reprinted by Chelsea, third edition, 1984.
relax44. G.E. Andrews, R. Askey and R. Roy, Special functions, Encyclopedia of Mathematics and its applications 71, Cambridge Univ. Press (1999). |
warning/0506/math0506065.html | ar5iv | text | # Sobolev Inequalities for Differential Forms and 𝐿_{𝑞,𝑝}-cohomology.
## 1. Introduction
Let us start by stating a Sobolev type Inequality for differential forms on a compact manifold:
###### Theorem 1.1.
Let $`(M,g)`$ be a smooth $`n`$-dimensional compact Riemannian manifold, $`1kn`$ and $`p,q(1,\mathrm{})`$. Then there exists a constant $`C`$ such that for any differential form $`\theta `$ of degree $`k1`$ on $`M`$ with coefficients in $`L^q`$, we have
(1.1)
$$\underset{\zeta Z^{k1}}{inf}\theta \zeta _{L^q(M)}Cd\theta _{L^p(M)},$$
if and only if
(1.2)
$$\frac{1}{p}\frac{1}{q}\frac{1}{n}.$$
Here $`Z^{k1}`$ denotes the set of smooth closed $`(k1)`$-forms on $`M`$.
The differential $`d\theta `$ in the inequality above is to be understood in the sense of currents.
Note that condition (1.2) is equivalent to
(1.3)
$$pn\text{or}p<n\text{ and }qp^{}=\frac{np}{np}.$$
In the case of zero forms (i.e. $`k=1`$), this theorem can be deduced from the corresponding result for functions with compact support in $`^n`$ by a simple argument using a partition of unity. The case of differential forms of higher degree can be proved using more involved reasoning based on standard results from the Hodge–De Rham theory and $`L^p`$-elliptic estimates obtained in the 1950’ by various authors. We give a sketch of such a proof in the appendix of this paper.
In the case of a non compact manifold, the inequality (1.1) is still meaningful if the differential form $`\theta `$ belongs to $`L^q`$. Although the condition (1.2) is still necessary in the non compact case, it is no longer sufficient and additional conditions must be imposed on the geometry of the manifold $`(M,g)`$ for a Sobolev inequality to hold.
The main goal of this paper is to investigate these conditions. Our Theorem 6.2 below gives a necessary and sufficient condition based on an invariant called the *$`L_{q,p}`$–cohomology of $`(M,g)`$* and which is defined as
$$H_{q,p}^k(M)=Z_p^k(M)/d\mathrm{\Omega }_{q,p}^{k1}(M).$$
where $`Z_p^k(M)`$ is the Banach space of closed $`k`$-forms $`\theta `$ in $`L^p(M)`$ and $`\mathrm{\Omega }_{q,p}^{k1}(M)`$ is the space of all $`(k1)`$-forms $`\varphi `$ in $`L^q(M)`$ such that $`d\varphi L^p`$.
We will also prove a regularization theorem saying that any $`L_{q,p}`$-cohomology class can be represented by a smooth form, provided that (1.2) holds (see Theorem 12.7). This implies in particular that the $`L_{q,p}`$-cohomology of a compact manifold $`M`$ coincides with the usual De Rham cohomology $`M`$ and it gives us a new proof of Theorem 1.1 above. This new proof is perhaps simpler than the classical one sketched in the appendix (at least it does not rely on the rather deep elliptic estimate).
The techniques of this paper also provide a proof of the following result which is a complement to Theorem 1.1:
###### Theorem 1.2.
Let $`(M,g)`$ be a smooth compact Riemannian manifold of dimension $`n`$ and $`p,q(1,\mathrm{})`$. There exists a constant $`C`$ such that for all closed differential forms $`\omega `$ of degree $`k`$ with coefficients in $`L^p(M)`$, there exists a differential form $`\theta `$ of degree $`k1`$ such that $`d\theta =\omega `$ and
(1.4)
$$\theta _{L^q}C\omega _{L^p},$$
if and only if $`p,q`$ satisfy the condition (1.2) and $`H_{\text{DeRham}}^k(M)=0`$.
Both Theorems 1.1 and 1.2 are proved at the end of section 12. In the non compact case, we prove in Theorem 6.1 below that the inequality (1.4) holds if and only if $`H_{q,p}^k(M,g)=0`$.
The Sobolev inequality is important because it is a key ingredient in solving partial differential equations. To illustrate this point, we show in section 13 how Theorem 6.2 can be used to solve the non linear equation
(1.5)
$$\delta (d\theta ^{p2}d\theta )=\alpha $$
for differential forms. Here $`\delta `$ is the formal adjoint to the exterior differential $`d`$.
Although it is certainly a nice observation that such Sobolev type inequalities for differential forms have interpretations in $`L_{q,p}`$-cohomology, this will not lead us very far unless we are able to compute some of this cohomology. Unfortunately, this is not an easy task and only few examples of $`L_{q,p}`$-cohomology groups are presently known. It is thus also one of our goals in this paper to begin developing some of the basic facts from the theory. In particular, we present here some results in the direction of duality (see section 8), a proof of the Poincaré Lemma for $`L_{q,p}`$-cohomology and a non vanishing result for the $`L_{q,p}`$-cohomology of the hyperbolic plane $`^2`$. This non vanishing result says in particular that the Sobolev inequality (1.4) for one-forms never holds on $`^2`$ for any $`p,q(1,\mathrm{})`$.
1.
2. Introduction
3. Definitions
4. Some elementary properties of $`L_{q,p}`$-cohomology
5. Banach complexes
6. $`L_{q,p}`$-cohomology and Banach complexes
7. $`L_{q,p}`$-cohomology and Sobolev inequality
8. Manifolds with finite volume and monotonicity
9. Almost duality
10. The $`L_{q,p}`$-cohomology of the line
11. The cohomology of the hyperbolic plane
12. The cohomology of the ball
13. Regularization of forms and cohomology classes
14. Relation with a non linear PDE
15. Torsion in $`L_2`$-cohomology and the Hodge-Kodaira decomposition
16. A “classic” proof of Theorem 1.1 in the compact case.
Let us shortly describe what is contained in the paper. In sections 2 and 3, we give the necessary definitions and we prove some elementary properties of $`L_{q,p}`$-cohomology. Then we present some basic facts of the theory of Banach complexes and we derive the cohomological interpretation of Sobolev inequalities for differential forms (section 4,5 and 6). In section 7, we prove some monotonicity properties for the $`L_{q,p}`$-cohomology of finite dimensional manifolds and in section 8 we introduce a notion of “almost duality” techniques (a standard Poincaré duality holds only when $`p=q`$). We apply these techniques to compute the $`L_{q,p}`$-cohomology of the line (section 9) and the hyperbolic plane (section 10) and to prove a version of the Poincaré Lemma (section 11). In section 12, we show that the $`L_{q,p}`$-cohomology of a manifold can be represented by smooth forms under the condition (1.2). Finally, we show in section 13 how the $`L_{q,p}`$-cohomology can be relevant in the study of some non linear PDE, and in section 14 we give a relation between the $`L_2`$-cohomology and the Laplacian on complete manifolds. The paper ends with an appendix describing an alternative proof of Theorems 1.1 based on $`L^p`$ elliptic estimates.
Remark. The reader might prefer to call the inequality (1.1) a *Poincaré inequality* and use the term *Sobolev inequality* only for the inequality (1.4). In fact there are various uses of the terms Poincaré and Sobolev inequalities. According to , the Poincaré inequality is simply a special case of the Sobolev one (it is in fact the case $`p=q`$). In this paper, we avoid the name Poincaré inequality.
*Acknowledgment.* Part of this research has been done in the autumn of 2001, when both authors stayed at IHES in Bures-Sur-Yvette. We are happy to thank the Institute for its warm hospitality. We also thank Pierre Pansu for his interest in our work and for the kindness and patience with which he explained us his viewpoint on the subject.
## 2. Definitions
Let us recall the notion of weak exterior differential of a differential form on a Riemannian manifold $`(M,g)`$.
We denote by $`C_c^{\mathrm{}}(M,\mathrm{\Lambda }^k)`$ the vector space of smooth differential forms of degree $`k`$ with compact support on $`M`$ and by $`L_{loc}^1(M,\mathrm{\Lambda }^k)`$ the space of differential $`k`$-forms whose coefficients (in any local coordinate system) are locally integrable.
###### Definition 2.1.
One says that a form $`\theta L_{loc}^1(M,\mathrm{\Lambda }^k)`$ is the *weak exterior differential* of a form $`\varphi L_{loc}^1(M,\mathrm{\Lambda }^{k1})`$ and one writes $`d\varphi =\theta `$ if for each $`\omega C_c^{\mathrm{}}(M,\mathrm{\Lambda }^{nk})`$, one has
$$_M\theta \omega =(1)^k_M\varphi d\omega .$$
Clearly $`d\varphi `$ is uniquely determined up to sets of Lebesgue measure zero, because $`d\varphi `$ is the exterior differential (in the sense of currents) of the current $`\varphi `$. It is also clear that $`dd=0`$, and this fact allows us to define various cohomology groups.
Let $`L^p(M,\mathrm{\Lambda }^k)`$ be the space of differential forms in $`L_{loc}^1(M,\mathrm{\Lambda }^k)`$ such that
$$\theta _p:=\left(_M|\theta |^p𝑑x\right)^{\frac{1}{p}}<\mathrm{}.$$
We then set $`Z_p^k(M):=L^p(M,\mathrm{\Lambda }^k)\mathrm{ker}d`$ (= the set of weakly closed forms in $`L^p(M,\mathrm{\Lambda }^k)`$) and
$$B_{q,p}^k(M):=d\left(L^q(M,\mathrm{\Lambda }^{k1})\right)L^p(M,\mathrm{\Lambda }^k).$$
###### Lemma 2.2.
$`Z_p^k(M)L^p(M,\mathrm{\Lambda }^k)`$ is a closed linear subspace. In particular it is a Banach space.
Proof We need to show that an arbitrary element $`\varphi \overline{Z}_p^k(M)`$ in the closure of $`Z_p^k(M)`$ is a weakly closed form. Choose a sequence $`\varphi _iZ_p^k(M)`$ such that $`\varphi _i\varphi `$ in $`L^p`$-norm. Since $`\varphi _i`$ are weakly closed forms, we have
$$_M\varphi _id\omega =0,$$
for any smooth differential forms $`\omega `$ of degree $`nk1`$ with compact support on $`M`$. Using Hölder’s inequality, we obtain
$$_M\varphi d\omega =_M(\varphi \varphi _i)d\omega \varphi \varphi _i_{L^p(M)}d\omega _{L^p^{}(M)}0.$$
Here $`1/p+1/p^{}=1`$.
Thus $`_M\varphi d\omega =0`$ for any $`\omega =C_c^{\mathrm{}}(M,\mathrm{\Lambda }^{nk1})`$ and hence $`\varphi Z_p^k(M)`$.
Observe that $`B_{q,p}^k(M)Z_p^k(M)`$ (because $`dd=0`$), we thus have
$$B_{q,p}^k(M)\overline{B}_{q,p}^k(M)Z_p^k(M)=\overline{Z}_p^k(M)L^p(M,\mathrm{\Lambda }^k).$$
###### Definition 2.3.
The $`L_{q,p}`$*-cohomology* of $`(M,g)`$ (where $`1p,q\mathrm{}`$) is defined to be the quotient
$$H_{q,p}^k(M):=Z_p^k(M)/B_{q,p}^k(M),$$
and the *reduced $`L_{q,p}`$-cohomology* of $`(M,g)`$ is
$$\overline{H}_{q,p}^k(M):=Z_p^k(M)/\overline{B}_{q,p}^k(M),$$
(where $`\overline{B}_{q,p}^k(M)`$ is the closure of $`B_{q,p}^k(M)`$). We also define the *torsion* as
$$T_{q,p}^k(M):=\overline{B}_{q,p}^k(M)/B_{q,p}^k(M).$$
We thus have the exact sequence
$$0T_{q,p}^k(M)H_{q,p}^k(M)\overline{H}_{q,p}^k(M)0.$$
The reduced cohomology is naturally a Banach space. The unreduced cohomology is a Banach space if and only if the torsion vanishes.
By Lemma 4.4 below, we see that the torsion $`T_{q,p}^k(M)`$ can be either $`\{0\}`$ or infinite dimensional. Indeed, if $`dimT_{q,p}^k(M)<\mathrm{}`$ then $`B_{q,p}^k(M)`$ is closed, hence $`T_{q,p}^k(M)=\{0\}`$. In particular, if $`dimT_{q,p}^k(M)0`$ then $`dimH_{q,p}^k(M)=\mathrm{}`$.
When $`p=q`$, we simply speak of $`L_p`$-cohomology and write $`H_p^k(M)`$ and $`\overline{H}_p^k(M)`$.
Example The $`L_{q,p}`$-cohomology of the bounded interval $`M=(0,1)`$ is easily computed: we clearly have $`H_{q,p}^0((0,1))=`$ and $`H_{q,p}^1((0,1))=0`$ for any $`1q,p\mathrm{}`$.
Indeed if $`\omega =a(x)dx`$ belongs to $`L^p((0,1))L^1((0,1))`$, then $`f(x):=_{\mathrm{}}^xa(s)𝑑s`$ belongs to $`L^q((0,1))`$ for any $`1q\mathrm{}`$.
The $`L_{q,p}`$-cohomology of the unbounded intervals and other examples will be computed below.
## 3. Some elementary properties of $`L_{q,p}`$-cohomology
### 3.1. Zero dimensional cohomology.
We have $`H_{q,p}^0(M)=\overline{H}_{q,p}^0(M)=Z_p^0(M)=H_p^0(M)`$ and these spaces have the following interpretation: $`dimH_{\mathrm{}}^0(M)`$ is the number of connected components of $`M`$ and $`dimH_p^0(M)`$ is the number of connected components with finite volume of $`M`$ if $`1p<\mathrm{}`$.
### 3.2. Conformal invariance.
Let $`(M,g)`$ be a Riemannian manifold of dimension $`n`$. Recall that a new metric $`g_1`$ is a conformal deformation of $`g`$ if $`g_1:=\rho ^2g`$ where $`\rho :M_+`$ is a smooth function.
The pointwise norms of a $`k`$-form $`\omega `$ with respect to the metrics $`g_1`$ and $`g`$ are related by the identity $`|\omega |_{g_1}=\rho ^k|\omega |_g`$. The volume elements are related by $`d\mathrm{vol}_{g_1}=\rho ^nd\mathrm{vol}_g`$. In particular
$$|\omega |_{g_1}^pd\mathrm{vol}_{g_1}=\rho ^{npk}|\omega |_g^pd\mathrm{vol}_g$$
for any $`k`$-form; likewise, $`|\theta |_{g_1}^qd\mathrm{vol}_{g_1}=\rho ^{nq(k1)}|\theta |_g^qd\mathrm{vol}_g`$ for any $`k1`$-form $`\theta `$. It follows that $`H_{q,p}^k(M,g_1)=H_{q,p}^k(M,g)`$ if $`npk=nq(k1)=0`$.
We thus have the
###### Theorem 3.1.
If $`q=\frac{n}{k1}`$ and $`p=\frac{n}{k}`$, then $`H_{q,p}^k(M,g)`$ and $`\overline{H}_{q,p}^k(M,g)`$ are conformal invariants.
## 4. Banach complexes
The abstract theory of Banach complexes is based on a combination of techniques from homological algebra and functional analysis; this theory is the natural framework of $`L_{q,p}`$-cohomology and we shall take this point of view to show the connections between Sobolev inequalities and $`L_{q,p}`$-cohomology.
There is not much literature on Banach complexes, we therefore give below all necessary definitions. The reader may look in for more information.
### 4.1. Cohomology of Banach complexes and abstract Sobolev inequalities.
###### Definition 4.1.
A *Banach complex* is a sequence $`F^{}=\{F^k,d_k\}_k`$ where $`F^k`$ is a Banach space, $`d_k:F^kF^{k+1}`$ is a bounded operator and $`d_{k+1}d_k=0`$.
Remarks 1.) It would be more correct to call such an object a Banach cocomplex (and to use the name complex for the case where $`d_k:F^kF^{k1}`$), but for simplicity, we shall speak of complexes.
2) To simplify notations, we usually note $`d`$ for any of the operators $`d_k`$.
###### Definition 4.2.
Given a Banach complex $`\{F^k,d\}`$ we introduce the following vector spaces:
* $`Z^k:=\mathrm{ker}(d:F^kF^{k+1})`$, it is a closed subspace of $`F^k`$;
* $`B^k:=`$Im$`(d:F^{k1}F^k)Z^k`$;
* $`H^k(F^{}):=Z^k/B^k`$ is the *cohomology* of the complex $`F^{}=\{F^k,d\}`$;
* $`\overline{H}^k(F^{}):=Z^k/\overline{B}^k`$ is *the reduced cohomology* of the complex $`F^{}`$;
* $`T^k(F^{}):=\overline{B}^k/B^k=H^k/\overline{H}^k`$ is the *torsion* of the complex $`F^{}`$.
Let us make a few elementary observations :
1. $`\overline{H}^k,Z^k`$ and $`\overline{B}^k`$ are Banach spaces;
2. The natural (quotient) topology on $`T^k:=\overline{B}^k/B^k`$ is coarse (any closed set is either empty or $`T^k`$);
3. We have the exact sequence
$$0T^kH^k\overline{H}^k0.$$
There is a natural notion of subcomplex:
###### Definition 4.3.
A *subcomplex* $`G^{}`$ of a Banach complex $`\{F^{},d\}`$ is a sequence of linear subspaces $`G^kF^k`$ (not necessarily closed) such that $`d(G^k)G^{k+1}`$. If all $`G^k`$ are closed subspaces, we say that $`G^{}`$ is a *Banach-subcomplex* of $`F^{}`$.
The cohomology of the subcomplex $`G^{}`$ is defined as
$$H^k(G^{})=(G^k\mathrm{ker}d)/d(G^{k1}).$$
Observe that in general $`H^k(G^{})`$ is not a Banach space, but there is no way to define a reduced cohomology of $`G^{}`$, unless $`G^{}F^{}`$ is a Banach-subcomplex.
###### Lemma 4.4.
For any Banach complex $`\{F^k,d\}`$, the following conditions are equivalent
1. $`T^k=0`$;
2. $`dimT_k<\mathrm{}`$;
3. $`B^kF^k`$ is closed.
Proof (i)$``$(ii) is obvious and (ii)$``$(iii) follows e.g from \[4, Th. 3.2 page 27\]. The implication (iii)$``$(i) follows directly from the definition of the torsion.
###### Proposition 4.5.
The following are equivalent:
1. $`H^k=0`$;
2. The operator $`d_{k1}:F^{k1}/Z^{k1}Z^k`$ admits a bounded inverse $`d_{k1}^1`$;
3. There exists a constant $`C_k`$ such that for any $`\theta Z^k`$ there is an element $`\eta F^{k1}`$ with $`d\eta =\theta `$ and
$$\eta _{F^{k1}}C_k\theta _{F^k}.$$
Proof (i) $``$ (ii). Suppose $`H^k=0`$. Then $`d_{k1}:F^{k1}/Z^{k1}Z^k`$ is a bijective bounded linear operator and by the open mapping theorem, the inverse map
$$d_{k1}^1:Z^kF^{k1}/Z^{k1}$$
is also a bounded operator.
(ii) $``$ (iii). Let $`\gamma `$ be the norm of $`d_{k1}^1:Z^kF^{k1}/Z^{k1}`$, then for any $`\theta Z^k`$ we can find $`\xi F^{k1}`$ such that $`d_{k1}\xi =\theta `$. Furthermore
$$[\xi ]_{F^{k1}/Z^{k1}}=\underset{\zeta Z^{k1}}{inf}\xi \zeta _{F^{k1}}\gamma \theta _{F^k}.$$
In particular, there exists $`\zeta Z^{k1}`$ such that $`\xi \zeta _{F^{k1}}2\gamma \theta _{F^k}`$. Let us set $`\eta :=(\xi \zeta )`$, then $`d_{k1}\eta =\theta `$ and $`\eta _{F^{k1}}C_k\theta _{F^k}`$with $`C_k=2\gamma =2d_{k1}^1_{Z^kF^{k1}/Z^{k1}}`$.
The implication (iii) $``$ (i) is clear.
###### Proposition 4.6.
The following conditions are equivalent:
1. $`T^k=0`$;
2. The operator $`d_{k1}:F^{k1}/Z^{k1}B^k`$ admits a bounded inverse $`d_{k1}^1`$.
And any one of these conditions imply
1. There exists a constant $`C_k^{^{}}`$ such that for any $`\xi F^{k1}`$ there is an element $`\zeta Z^{k1}`$ such that
(4.1)
$$\xi \zeta _{F^{k1}}C_k^{^{}}d\xi _{F^k}.$$
Proof The conditions (i) and (ii) are equivalent, because the existence of a bounded inverse operator is equivalent to the closedness of $`B^{k1}`$ by the open mapping theorem.
Let us assume that $`T^k=0`$ and prove (iii). By hypothesis, $`B^k`$ is a Banach space and $`d_{k1}:F^{k1}/Z^{k1}B^k`$ is a bijective bounded linear operator. Thus, by the open mapping theorem, the inverse $`d_{k1}^1:B^kF^{k1}/Z^{k1}`$ is also a bounded operator.
Let $`\gamma `$ be the norm of $`d_{k1}^1:B^kF^{k1}/Z^{k1}`$, then for any $`\xi F^{k1}`$ we have
$$[\xi ]_{F^{k1}/Z^{k1}}=\underset{\zeta Z^{k1}}{inf}\xi \zeta _{F^{k1}}\gamma d_{k1}\xi _{F^k}$$
in particular, there exists $`\zeta Z^{k1}`$ such that $`\xi \zeta _{F^{k1}}2\gamma d_{k1}\xi _{F^k}`$.
###### Proposition 4.7.
If $`F^{k1}`$ is a reflexive Banach space, then the three conditions of the previous proposition are equivalent.
Proof We only need to show that (iii)$``$ (i) i.e. $`B^k=\overline{B}^kF^k`$ provided (4.1) holds and $`F^{k1}`$ is a reflexive. Let $`\theta \overline{B}^k`$, then there exists a sequence $`\xi _iF^{k1}`$ such that $`d_{k1}\xi _i\theta `$ in $`F^k`$. By hypothesis there exists a sequence $`\zeta _iZ^{k1}`$ such that $`\xi _i\zeta _i_{F^{k1}}C_k^{}d\xi _i_{F^k}`$. In particular, the sequence $`\{\eta _i:=(\xi _i\zeta _i)\}`$ is bounded, we may thus find a subsequence (still denoted $`\{\eta _i\}`$) which converges weakly to an element $`\eta F^{k1}`$.
Using the Mazur Lemma (see e.g. chap. V §1, Theorem 2, page 120 in ), we may construct a sequence $`\{\stackrel{~}{\eta }_i=_{j=i}^{N(i)}a_i\eta _j\}`$ of convex combinations of $`\eta _i`$ such that $`\stackrel{~}{\eta }_i`$ converges strongly to $`\eta `$. We then have
$$d_{k1}\eta =\underset{i\mathrm{}}{lim}d_{k1}\stackrel{~}{\eta }_i=\underset{i\mathrm{}}{lim}\underset{j=i}{\overset{N(i)}{}}a_id_{k1}\eta _i=\underset{i\mathrm{}}{lim}\underset{j=i}{\overset{N(i)}{}}a_id_{k1}\xi _j=\theta $$
hence $`\theta \text{Im}(d)=B^k`$. We proved that $`B^k`$ is closed, i.e. $`T^k=0`$.
### 4.2. Morphisms and homotopies of Banach complexes.
This part will be useful to regularize $`L_{q,p}`$-cohomology, see section 12.
Definitions 1) A *morphism* $`R^{}`$ between two Banach complexes $`F^{}=\{F^k,d\}`$ and $`E^{}=\{E^k,d\}`$ is a family of bounded operators $`R^k:F^kE^k`$ such that
$$d_kR^k=R^{k+1}d_k.$$
2) A *homotopy* between two morphisms $`R^{}`$ and $`S^{}:F^{}E^{}`$ is a family of bounded operators $`A^k:F^kE^{k1}`$ such that
$$S^kR^k=d_{k1}A^k+A^{k+1}d_k.$$
3) A *weak homotopy* between two morphisms $`R^{}`$ and $`S^{}:F^{}E^{}`$ is a sequence of families of bounded operators $`A_j^k:F^kE^{k1}`$ such that for any element $`xF^k`$ we have
$$\underset{j\mathrm{}}{lim}(d_{k1}A_j^k+A_j^{k+1}d_k)x(S^kR^k)x=0.$$
Observe that, if $`R^{}=\{R^k:F^kE^k\}`$ is a morphism, then its image is a subcomplex of $`E^{}`$ and it is a Banach-subcomplex if and only if all $`R^k`$ are closed operators. The kernel of $`R^{}`$ is always a Banach-subcomplex of $`F^{}`$.
###### Proposition 4.8.
Let $`R^{}:F^{}F^{}`$ be an endomorphism of a Banach complex $`\{F^{},d\}`$ such $`R^{}(F^{})G^{}`$ where $`G^{}`$ is a subcomplex.
If there exists a homotopy $`\{A^k:F^kF^{k1}\}`$ between $`R^{}`$ and the identity operator $`I:F^{}F^{}`$, then
$$H^k(F^{})=H^k(G^{}).$$
Proof Given $`\xi Z^k(F^{})`$, we observe that $`R^k\xi Z^k(G^{})`$ because $`dR\xi =Rd\xi =0`$. If $`\xi =d\eta B^k(F^{})`$, then $`R^k\xi =R^kd\eta =dR^k\eta B^k(G^{})`$.
This proves that $`[R\xi ]`$ is a well defined cohomology class in $`H^k(G^{})`$ for any cohomology class $`[\xi ]H^k(F^{})`$.
But since
$$\xi R\xi =dA\xi +Ad\xi =dA\xi $$
for any $`\xi Z^k(F^{})`$, we see that in fact $`[R\xi ]=[\xi ]H^k(F^{})`$ and the Proposition is proved.
The following result is a generalization of the previous proposition.
###### Proposition 4.9.
(1) Any morphism $`R^{}:F^{}E^{}`$ between two Banach complexes induces a sequence of linear homomorphisms $`H^kR^{}:H^k(F^{})H^k(E^{})`$ from the cohomology of $`F^{}`$ to the cohomology of $`E^{}`$.
(2) The morphism $`R^{}:F^{}E^{}`$ induces a sequence of bounded operators $`\overline{H}^kR^{}:\overline{H}^k(F^{})\overline{H}^k(E^{})`$ from the reduced cohomology of $`F^{}`$ to the reduced cohomology of $`E^{}`$.
(3) If there exists a homotopy between two morphisms $`R^{}`$ and $`S^{}:F^{}E^{}`$, then the corresponding homomorphisms on the cohomology groups coincide:
$$H^kR^{}=H^kS^{}:H^k(F^{})H^k(E^{}).$$
(4) If there exists a weak homotopy between two morphisms $`R^{}`$ and $`S^{}:F^{}E^{}`$, then the corresponding morphisms on the reduced cohomology groups coincide:
$$\overline{H}^kR^{}=\overline{H}^kS^{}:\overline{H}^k(F^{})\overline{H}^k(E^{}).$$
Proof (1) Because $`dR^{}=R^{}d`$, the image $`R^{}([\omega ])`$ of any cohomology class $`[\omega ]`$ of the complex $`F^{}`$ is a well defined cohomology class of the complex $`E^{}`$.
(2) Using the continuity of $`R^{}`$ and $`dR^{}=R^{}d`$, we see that closure of the image $`R^{}([\omega ])`$ of a reduced cohomology class of $`F^{}`$ is a well defined reduced cohomology class of $`E^{}`$. By the boundedness of $`R^k`$, the operators $`\overline{H}^kR^{}:\overline{H}^k(F^{})\overline{H}^k(E^{})`$ is also bounded.
(3) The condition $`S^kR^k=dA^k+A^{k+1}d`$ implies that for any $`\xi Z^k(F^{})`$ we have $`\left(S^k\xi R^k\xi \right)=d(A^k\xi )B^k(E^{})`$.
(4) The condition $`lim_j\mathrm{}(dA_j^k+A_j^{k+1}d)x(S^kR^k)x=0`$ for any $`xF^k`$ implies that for any $`\xi Z^k(F^{})`$ we have
$$\underset{j\mathrm{}}{lim}S^k\xi R^k\xi d(A_j^k\xi )=0.$$
A special case of the previous Proposition is given in the following definitions:
###### Definition 4.10.
a) A Banach complex $`\mathrm{F}^{}=\{\mathrm{F}^\mathrm{k},\mathrm{d}\}`$ is *acyclic* if there exists a family of bounded operators $`\mathrm{A}^\mathrm{k}:\mathrm{F}^\mathrm{k}\mathrm{F}^{\mathrm{k}1}`$ such that
$$Id=dA^k+A^{k+1}d.$$
b) The Banach complex $`\mathrm{F}^{}`$ is *weakly acyclic* if for any $`\mathrm{k}`$ there exists a sequence of bounded operators $`\mathrm{A}_\mathrm{j}^\mathrm{k}:\mathrm{F}^\mathrm{k}\mathrm{F}^{\mathrm{k}1}`$ such that for any element $`\mathrm{x}\mathrm{F}^\mathrm{k}`$ we have
$$\underset{j\mathrm{}}{lim}(dA_j^k+A_j^{k+1}d)xx=0.$$
In other words, $`F^{}`$ is (weakly) acyclic if and only if there exists a (weak) homotopy from the identity $`Id:F^{}F^{}`$ to the trivial morphism $`0:F^{}F^{}`$ It is thus clear that an acyclic complex has trivial cohomology and a weakly acyclic complex has trivial reduced cohomology.
## 5. $`L_{q,p}`$-cohomology and Banach complexes
In this section, we explain how the $`L_{q,p}`$-cohomology of a Riemannian manifold $`(M,g)`$ can be formally seen as the cohomology of some complex of Banach spaces. Let us start by introducing the notation
$$\mathrm{\Omega }_{q,p}^k(M):=\left\{\omega L^q(M,\mathrm{\Lambda }^k)\right|d\omega L^p\}.$$
This is a Banach space for the graph norm
(5.1)
$$\omega _{\mathrm{\Omega }_{q,p}}:=\omega _{L^q}+d\omega _{L^p}.$$
By standard arguments of functional analysis (see e.g. ) , it can be proved that $`\mathrm{\Omega }_{q,p}^k(M)`$ is a reflexive Banach space for any $`1<p,q<\mathrm{}`$. We will also prove in section 12 that smooth forms are dense in $`\mathrm{\Omega }_{q,p}^k(M)`$ for any $`1p,q<\mathrm{}`$.
To define a Banach complex, we choose an arbitrary finite sequence of numbers
$$\pi =\{p_0,p_1,\mathrm{},p_n\}[1,\mathrm{}],$$
and define
$$\mathrm{\Omega }_\pi ^k(M):=\mathrm{\Omega }_{p_k,p_{k+1}}^k(M).$$
Observe that $`\mathrm{\Omega }_\pi ^n(M)=L^{p_n}(M,\mathrm{\Lambda }^n)`$ and $`\mathrm{\Omega }_{p,p}^1(M)`$ coincides with the Sobolev space $`W^{1,p}(M)`$.
Since the exterior differential is a bounded operator $`d:\mathrm{\Omega }_\pi ^{k1}\mathrm{\Omega }_\pi ^k`$, we have constructed a Banach complex.
$$0\mathrm{\Omega }_\pi ^0\stackrel{𝑑}{}\mathrm{}\stackrel{𝑑}{}\mathrm{\Omega }_\pi ^{k1}\stackrel{𝑑}{}\mathrm{\Omega }_\pi ^k\stackrel{𝑑}{}\mathrm{}\stackrel{𝑑}{}\mathrm{\Omega }_\pi ^n0.$$
###### Definition 5.1.
*The (reduced) $`L_\pi `$-cohomology of $`M`$ is the (reduced) cohomology of the Banach complex* $`\{\mathrm{\Omega }_\pi ^k(M),d_k\}`$*.*
The $`L_\pi `$-cohomology space $`H_\pi ^k(M)`$ depends only on $`p_{k\text{ }}`$and $`p_{k1}`$ and we have in fact
$$H_\pi ^k(M)=H_{p_{k1},p_k}^k(M)\text{ and }\overline{H}_\pi ^k(M)=\overline{H}_{p_{k1},p_k}^k(M).$$
Two cases are of special interest:
1. The $`L_p`$-cohomology, which corresponds to the constant sequence $`\pi =\{p,p,\mathrm{},p\}`$.
2. The *conformal cohomology*, which corresponds to the sequence $`p_0=\mathrm{}`$, and $`p_k=\frac{n}{k}`$ for $`k=1,\mathrm{},n`$. The cohomology associated to this sequence is a conformal invariant of the manifold by Theorem 3.1.
Let us remark here that $`\left(\frac{1}{p_k}\frac{1}{p_{k1}}\right)=\frac{1}{n}`$.
## 6. $`L_{q,p}`$-cohomology and Sobolev inequality
We are now in position to give the interpretation of $`L_{q,p}`$-cohomology in terms of a Sobolev type inequality for differential forms on a Riemannian manifold $`(M,g)`$:
###### Theorem 6.1.
$`H_{q,p}^k(M,g)=0`$ if and only if there exists a constant $`C<\mathrm{}`$ such that for any closed $`p`$-integrable differential form $`\omega `$ of degree $`k`$ there exists a differential form $`\theta `$ of degree $`k1`$ such that $`d\theta =\omega `$ and
$$\theta _{L^q}C\omega _{L^p}.$$
This result is a direct consequence of Proposition 4.5.
###### Theorem 6.2.
A) If $`T_{q,p}^k(M)=0`$, then there exists a constant $`C^{}`$ such that for any differential form $`\theta \mathrm{\Omega }_{q,p}^{k1}(M)`$ of degree $`k1`$ there exists a closed form $`\zeta Z_q^{k1}(M)`$ such that
(6.1)
$$\theta \zeta _{L^q}C^{}d\theta _{L^p}.$$
B) Conversely, if $`1<q<\mathrm{}`$, and if there exists a constant $`C^{}`$ such that for any form $`\theta \mathrm{\Omega }_{q,p}^{k1}(M)`$ of degree $`k1`$ there exists $`\zeta Z_q^{k1}(M)`$ such that (6.1) holds, then $`T_{q,p}^k(M)=0`$.
This statement follows immediately from Proposition 4.6 and 4.7.
## 7. Manifolds with finite volume and monotonicity
The $`L_{q,p}`$-cohomology of a manifold with finite volume has some monotonicity properties. In the next statement, the symbol $`H_2H_1`$ (where $`H_1,H_2`$ are vector spaces) means that $`H_1`$ is a quotient of $`H_2`$.
###### Proposition 7.1.
If $`(M,g)`$ has finite volume, $`1p\mathrm{}`$ and $`1q_1q_2\mathrm{}`$, then $`\overline{H}_{q_2,p}^k(M)\overline{H}_{q_1,p}^k(M)`$ and $`H_{q_2,p}^k(M)H_{q_1,p}^k(M)`$.
Proof Since $`1q_1q_2`$ and $`M`$ has finite volume, we have $`L^{q_1}(M,\mathrm{\Lambda }^k)L^{q_2}(M,\mathrm{\Lambda }^k)`$, hence $`\mathrm{\Omega }_{q_1,p}^{k1}\mathrm{\Omega }_{q_2,p}^{k1}`$ and thus
$`\overline{B}_{q_1,p}^k(M)`$ $`=`$ $`\overline{d\left(\mathrm{\Omega }_{q_1,p}^{k1}\right)}L^p(M,\mathrm{\Lambda }^k)`$
$``$ $`\overline{d\left(\mathrm{\Omega }_{q_2,p}^{k1}\right)}L^p(M,\mathrm{\Lambda }^k)`$
$`=`$ $`\overline{B}_{q_2,p}^k(M).`$
Since $`B_2B_1Z`$ implies $`Z/B_1Z/B_2`$, we have
$$\overline{H}_{q_2,p}^k(M)=Z_p^k/\overline{B}_{q_2,p}^k(M)Z_p^k/\overline{B}_{q_1,p}^k(M)=\overline{H}_{q_1,p}^k(M).$$
The proof for unreduced cohomology is the same.
We also have some kind of monotonicity with respect to $`p`$:
###### Proposition 7.2.
If $`(M,g)`$ has finite volume $`1p_2p_1\mathrm{}`$ and $`1q_1q_2\mathrm{}`$, then
$$H_{q_2,p_2}^k(M)=0H_{q_1,p_1}^k(M)=0.$$
Proof Since $`M`$ has finite volume, $`q_1q_2`$ and $`p_2p_1`$, we have<sup>1</sup><sup>1</sup>1The symbol $``$ means that the inequality holds up to some constant. for any $`q_2`$-integrable form $`\theta `$ and any $`p_1`$-integrable form $`\omega `$
$$\theta _{L^{q_1}}\theta _{L^{q_2}}\text{ and }\omega _{L^{p_2}}\omega _{L^{p_1}}.$$
Since $`H_{q_2,p_2}^k(M)=0`$, we know from Theorem 6.1 that for any closed $`p_2`$-integrable form $`\omega `$ of degree $`k`$ there exists a differential form $`\theta `$ of degree $`k1`$ such that $`d\theta =\omega `$ and
$$\theta _{L^{q_2}}\omega _{L^{p_2}}.$$
Combining this inequality with two previous inequalities we get
$$\theta _{L^{q_1}}\omega _{L^{p_1}}$$
and the result immediately follows from the same Theorem 6.1.
For the torsion, we need to avoid the values $`q=1`$ and $`q=\mathrm{}`$:
###### Proposition 7.3.
If $`(M,g)`$ has finite volume $`1p_2p_1\mathrm{}`$ and $`1<q_1q_2<\mathrm{}`$, then
$$T_{q_2,p_2}^k(M)=0T_{q_1,p_1}^k(M)=0.$$
Proof Again, since $`q_1q_2`$ we have $`\zeta Z_{q_2}^{k1}(M)\zeta Z_{q_1}^{k1}(M)`$ and
$$\theta \zeta _{L^{q_1}}\theta \zeta _{L^{q_2}}\text{ and }d\theta _{L^{p_2}}d\theta _{L^{p_1}}.$$
We may thus argue as in the previous proof using Theorem 6.2.
## 8. Almost duality
It has been proved in that for complete manifolds the dual space of $`\overline{H}_p^k(M)`$ coincides with $`\overline{H}_p^{^{}}^{nk}(M)`$ where $`\frac{1}{p}+\frac{1}{p^{^{}}}=1`$ (there is also a duality result for non complete manifolds). The duality is based on the pairing $`_M\alpha \beta `$ where $`\alpha \mathrm{\Omega }_p^k(M)`$ and $`\beta \mathrm{\Omega }_p^{^{}}^k(M)`$.
For $`L_{q,p}`$-cohomology we have no convenient description of dual spaces, but the notion of *almost duality* which we now introduce is sufficient for many calculations.
We start with a rather elementary result about the non vanishing of $`L_{q,p}`$-cohomology:
###### Lemma 8.1.
Let $`(M,g)`$ be an arbitrary Riemannian manifold of dimension $`n`$. Let $`\alpha Z_p^k(M)`$. If there exists $`\gamma C_\text{c}^{\mathrm{}}(M,\mathrm{\Lambda }^{nk})`$ such that $`d\gamma =0`$ and $`_M\alpha \gamma 0`$, then $`[\alpha ]0`$ in $`\overline{H}_{q,p}^k(M)`$ for any $`1q\mathrm{}`$.
Proof Suppose that $`\alpha \overline{B}_{q,p}^k(M)`$. Then $`\alpha =\underset{j\mathrm{}}{lim}d\beta _j`$ (where the limit is in $`L^p`$-topology) for some $`\beta _jL^q(M,\mathrm{\Lambda }^{k1})`$ with $`d\beta _jL^p(M,\mathrm{\Lambda }^k)`$. We then have for any closed form with compact support $`\gamma C_\text{c}^{\mathrm{}}(M,\mathrm{\Lambda }^{nk})`$
$$_M\gamma \alpha =\underset{j\mathrm{}}{lim}_M\gamma d\beta _j=\underset{j\mathrm{}}{lim}(1)^{nk+1}_M𝑑\gamma \beta _j=0$$
in contradiction to the assumption. ∎
There are several generalizations of this result :
###### Proposition 8.2.
Let $`(M,g)`$ be an arbitrary Riemannian manifold of dimension $`n`$. Let $`\alpha Z_p^k(M)`$. Then
A) If there exists a sequence $`\{\gamma _i\}C_\text{c}^{\mathrm{}}(M,\mathrm{\Lambda }^{nk})`$ such that
1. $`\underset{i\mathrm{}}{lim\; inf}{\displaystyle _M}\alpha \gamma _i>0`$;
2. $`\underset{i\mathrm{}}{lim}d\gamma _i_q^{}=0`$ where $`q^{}=\frac{q}{q1}`$.
Then $`[\alpha ]0`$ in $`H_{q,p}^k(M)`$.
B) If there exists a sequence $`\{\mathrm{\gamma }_\mathrm{i}\}\mathrm{C}_\text{c}^{\mathrm{}}(\mathrm{M},\mathrm{\Lambda }^{\mathrm{n}\mathrm{k}})`$ satisfying the conditions (i) and (ii) above and
1. $`\gamma _i_p^{}`$ is a bounded sequence for $`p^{}=\frac{p}{p1}`$.
Then $`[\alpha ]0`$ in $`\overline{H}_{q,p}^k(M)`$.
Proof A) Suppose that $`\alpha =d\beta `$ for some $`\beta L^q(M,\mathrm{\Lambda }^{k1})`$, then by Hölder inequality we have for any $`\gamma C_\text{c}^{\mathrm{}}(M,\mathrm{\Lambda }^{nk})`$
$$\left|_M\alpha \gamma \right|=\left|_M𝑑\beta \gamma \right|=\left|_M\beta d\gamma \right|\beta _qd\gamma _q^{}.$$
It follows that for any sequence $`\{\gamma _i\}C_\text{c}^{\mathrm{}}(M,\mathrm{\Lambda }^{nk})`$ such that $`lim_i\mathrm{}d\gamma _q^{}=0`$, we have $`\underset{i\mathrm{}}{lim}\left|{\displaystyle _M}\alpha \gamma \right|\underset{i\mathrm{}}{lim}\beta _qd\gamma _i_{L^q^{}(M)}=0`$.
B) Suppose that $`\alpha \overline{B}_{q,p}^k(M)`$. Then $`\alpha =\underset{j\mathrm{}}{lim}d\beta _j`$ for $`\beta _jL^q(M,\mathrm{\Lambda }^{k1})`$ with $`d\beta _jL^p(M,\mathrm{\Lambda }^k)`$. We have for any $`i,j`$
$$_M\gamma _i\alpha =_M\gamma _id\beta _j+_M\gamma _i(\alpha d\beta _j).$$
For each $`j`$, we can find $`i=i(j)`$ large enough so that $`d\gamma _{i(j)}_q^{}\beta _j_q1/j`$, we thus have
$$|_M\gamma _{i(j)}d\beta _j||_Md\gamma _{i(j)}\beta _j|d\gamma _{i(j)}_q^{}\beta _j_q\frac{1}{j}.$$
On the other hand
$$\underset{j\mathrm{}}{lim}|_M\gamma _{i(j)}(\alpha d\beta _j)|\underset{j\mathrm{}}{lim}\gamma _{i(j)}_p^{}(\alpha d\beta _j)_p=0$$
since $`\gamma _{i(j)}_p^{}`$ is a bounded sequence and $`(\alpha d\beta _j)_p0`$. It follows that $`_M\gamma _{i(j)}\alpha 0`$ in contradiction to the hypothesis.
### 8.1. The case of complete manifolds
If $`M`$ is a complete manifold, we don’t need to assume that the form $`\gamma `$ from the previous discussion has compact support.
###### Proposition 8.3.
Assume that $`M`$ is complete. Let $`\alpha Z_p^k(M)`$, and assume that there exists a smooth closed $`(nk)`$-form $`\gamma `$ such that $`\gamma Z_q^{}^{nk}(M)`$, for $`q^{}=\frac{q}{q1},\gamma \alpha L^1(M)`$ and
$$_M\gamma \alpha 0,$$
then $`\alpha B_{q,p}^k(M)`$. In particular, $`H_{q,p}^k(M)\mathrm{}`$.
This proposition has also version for reduced $`L_{q,p}`$-cohomology:
###### Proposition 8.4.
Assume that $`M`$ is complete. Let $`\alpha Z_p^k(M)`$, and assume that there exists a smooth closed $`(nk)`$-form $`\gamma Z_p^{}^{nk}(M)Z_q^{}^{nk}(M)`$, where $`p^{}=\frac{p}{p1}`$ and $`q^{}=\frac{q}{q1}`$, such that
$$_M\gamma \alpha 0,$$
then $`\alpha \overline{B}_{q,p}^k(M)`$ where $`q^{}=\frac{q}{q1}`$. In particular, $`\overline{H}_{q,p}^k(M)\mathrm{}`$.
The proofs are based on the following integration by part lemma:
###### Lemma 8.5.
Assume that $`M`$ is complete. Let $`\beta L^q(M,\mathrm{\Lambda }^{k1})`$ be such that $`d\beta L^p(M,\mathrm{\Lambda }^k)`$, and $`\gamma L^p^{^{}}(M,\mathrm{\Lambda }^{nk})`$ be such that $`d\gamma L^q^{^{}}(M,\mathrm{\Lambda }^{nk+1})`$ where $`\frac{1}{p}+\frac{1}{p^{}}=\frac{1}{q}+\frac{1}{q^{}}=1`$.
If $`\gamma `$ is smooth and $`\gamma d\beta L^1(M)`$, then
(8.1)
$$_M\gamma d\beta =(1)^{nk+1}_M𝑑\gamma \beta ,$$
In particular, if $`\gamma L_p^{^{}}^{nk}(M)L_q^{^{}}^{nk+1}(M)`$, then the above conclusion holds.
Proof The integrability of $`d\gamma \beta `$ and $`\gamma d\beta `$ is a direct consequence of Hölder’s inequality.
By Hölder’s inequality, the forms $`d\gamma \beta `$ and $`\gamma d\beta `$ both belong to $`L^1(M)`$.
If $`\gamma `$ is a smooth form with compact support, then the equation (8.1) follows from the definition of the weak exterior differential (of $`\beta `$).
If the support of $`\gamma `$ is not compact, we set $`\gamma _i:=\psi _i\gamma `$ where $`\{\psi _i\}`$ is a sequence of smooth functions with compact support such that $`\psi _i(x)1`$ uniformly on every compact subset, $`0\psi _i(x)1`$ and $`|d\psi _i|_x1`$ for all $`xM`$ (such a sequence exists on any complete manifold).
The formula (8.1) holds for each $`\gamma _i`$ (since these forms have compact support).
Using $`|d\psi _i|_x1`$, we have the estimate
$$|\gamma _id\beta +(1)^{nk}d\gamma _i\beta ||d\gamma \beta |+|\gamma d\beta |+|\gamma \beta |L^1(M).$$
By Lebesgue’s dominated convergence theorem, we thus have
$$_M\left(\gamma d\beta +(1)^{nk}d\gamma \beta \right)=\underset{i\mathrm{}}{lim}_M\left(\gamma _id\beta +(1)^{nk}d\gamma _i\beta \right)=0.$$
Proof of Proposition 8.3 Suppose that $`\alpha B_{q,p}^k(M)`$. Then $`\alpha =d\beta `$ for some $`\beta L^q(M,\mathrm{\Lambda }^{k1})`$. By the previous lemma, we have
$$_M\gamma \alpha =_M\gamma d\beta =(1)^{nk+1}_M𝑑\gamma \beta =0$$
(since $`\gamma `$ is closed) in contradiction to the assumption. ∎
Proof of Proposition 8.4 Suppose that $`\alpha \overline{B}_{q,p}^k(M)`$. Then $`\alpha =\underset{j\mathrm{}}{lim}d\beta _j`$ (where the limit is in $`L^p`$-topology) for some $`\beta _jL^q(M,\mathrm{\Lambda }^{k1})`$ with $`d\beta _jL^p(M,\mathrm{\Lambda }^k)`$. Since $`d\gamma =0`$, we have
$$_M\gamma \alpha =\underset{j\mathrm{}}{lim}_M\gamma d\beta _j=\underset{j\mathrm{}}{lim}(1)^{nk+1}_M𝑑\gamma \beta _j=0,$$
which contradicts our hypothesis. ∎
## 9. The $`L_{q,p}`$-cohomology of the line
In the following three sections, we compute the $`L_{q,p}`$-cohomology of the line, the hyperbolic plane and the ball. We will see in particular that the only case where $`H_{q,p}^1()`$ vanishes is when $`q=\mathrm{}`$, $`p=1`$ :
###### Proposition 9.1.
$`H_{\mathrm{},1}^1()=0`$.
Proof If $`\omega =a(x)dx`$ belongs to $`L^1()`$, then $`f(x):=_{\mathrm{}}^xa(s)𝑑s`$ belongs to $`L^{\mathrm{}}()`$, hence $`H_{1,\mathrm{}}^1()=0`$.
###### Proposition 9.2.
$`T_{q,p}^1()0`$ for any $`1p,q\mathrm{}`$ with the only exception of $`q=\mathrm{}`$, $`p=1`$.
Proof Assume first that $`q<\mathrm{}`$. We know from Theorem 6.2 that if we had $`T_{q,p}^1()=0`$, then there would exist a Sobolev inequality for functions on the real line $``$:
(9.1)
$$\underset{z}{inf}\left(_{\mathrm{}}^{\mathrm{}}|f(x)z|^q𝑑x\right)^{1/q}C\left(_{\mathrm{}}^{\mathrm{}}|f^{}(x)|^p𝑑x\right)^{1/p}$$
for some constant $`C<\mathrm{}`$.
To see that no such inequality is possible, consider a family of smooth functions with compact support $`f_a:`$ such that $`f(x)=1`$ if $`x[1,a]`$ and $`f_a(x)=0`$ if $`x[0,a+1]`$. We may also assume that $`f_a^{}_L^{\mathrm{}}2`$. Assume now that the inequality (9.1) holds. Then the constant $`z`$ must be zero and we have
$$_{\mathrm{}}^{\mathrm{}}|f_a(x)|^q𝑑xa1\text{and}_{\mathrm{}}^{\mathrm{}}|f_a^{}(x)|^p𝑑x2^{1+p},$$
hence
$$C2^{1\frac{1}{p}}(a1)^{\frac{1}{q}}$$
for all $`a>0`$ and we conclude that $`C=\mathrm{}`$.
Assume now that $`q=\mathrm{}`$ and $`p>1`$. Again, if we had $`T_{\mathrm{},p}^1()=0`$, there would exist $`C<\mathrm{}`$ such that for any $`fL^p()`$:
(9.2)
$$\underset{z}{inf}f(x)z_{\mathrm{}}Cf^{}(x)_{L^p()}.$$
Let us consider the functions $`g_k(x):=e^{\pi kx^2}`$ and $`f(x):=_{\mathrm{}}^xg(u)𝑑u`$.
We have $`0f(x)<supf=_{\mathrm{}}^{\mathrm{}}g(u)𝑑u=\frac{1}{\sqrt{k}}`$, hence $`inf_zf(x)z_{\mathrm{}}=\frac{1}{2\sqrt{k}}`$. On the other hand $`f^{}(x)_{L^p()}=(kp)^{1/2p}`$, hence the constant in (9.2) satisfies
$$\frac{1}{2}k^{1/2}C(kp)^{1/2p}$$
for all $`k>0`$, i.e. $`C=\mathrm{}`$ since $`p>1`$.
Finally, we have $`T_{\mathrm{},1}^1()=0`$ since $`H_{\mathrm{},1}^1()=0`$.
Let us turn to the reduced cohomology:
###### Proposition 9.3.
$`\overline{H}_{q,p}^1()0`$ if and only if $`p=1`$ and $`1q<\mathrm{}`$.
Proof For $`p=1,q=\mathrm{}`$, we have $`\overline{H}_{\mathrm{},1}^1()=H_{\mathrm{},1}^1()=0`$.
Assume $`1q\mathrm{}`$ and $`1<p\mathrm{}`$ and let $`\omega =a(x)dxL^p()`$. For each $`m`$, we set $`\omega _m:=\chi _{[m,m]}\omega =(\chi _{[m,m]}(x)a(x))dx`$. Let us choose a continuous function $`\lambda _m(x)`$ with compact support in $`[0,\mathrm{})`$ such that $`_{}\lambda _m(x)𝑑x=_m^ma(x)𝑑x`$ and $`\lambda _m_{L^p()}<\frac{1}{m}`$.
Let $`b_m(x):=_{\mathrm{}}^x\left(\chi _{[m,m]}(t)a(t)\lambda _m(t)\right)𝑑t`$, then $`b_mL^q()`$ (in fact $`b_m`$ has compact support) and $`db_m\omega _{L^p()}a_{L^p([m,m])}+\lambda _m_{L^p()}0`$ as $`m\mathrm{}`$. This shows that $`\overline{H}_{q,p}^1()=0`$.
Assume now that $`p=1`$ and $`1q<\mathrm{}`$ and let $`\omega =a(x)dx`$ be a 1-form on $``$ such that $`_{}f\omega =1`$ and $`a(x)`$ is smooth with compact support (say $`\mathrm{supp}(a)[1,2]`$). Let $`f_j:`$ be a sequence of smooth functions with compact support such that $`f_j=1`$ on $`[1,2]`$, $`f_j_L^{\mathrm{}}=1`$ and $`f_j^{}_{L^q^{}}\frac{1}{j}`$ where $`q^{}=q/(q1)`$.
Using Proposition 8.2, we see that $`[\omega ]0\overline{H}_{q,1}^1()`$, because $`\omega L^1()`$ and the sequence $`\{f_j\}C_c^{\mathrm{}}()`$ satisfies the three conditions of that Proposition.
Remarks 1.) In degree 0, the $`L_{q,p}`$-cohomology is controlled by the volume: $`\overline{H}_{q,p}^0()=H_{q,p}^0()=0`$ if and only if $`p<\mathrm{}`$ and $`\overline{H}_{q,\mathrm{}}^0()=H_{q,\mathrm{}}^{\mathrm{}}()=`$.
2.) All the results of this section also hold for the half-line $`_+`$.
## 10. The cohomology of the hyperbolic plane
We treat in this section the case of the hyperbolic plane.
Recall that the hyperbolic plane is the Riemannian manifold $`^2=\{(u,v)^2:v>0\}`$ with the metric $`ds^2=v^2(du^2+dv^2)`$.
###### Theorem 10.1.
For any $`q,p(1,\mathrm{})`$ we have
$$dim(\overline{H}_{q,p}^1(^2))=\mathrm{}.$$
It will be convenient to introduce new coordinates (the so called “horocyclic coordinates”) $`y:=u`$, $`z:=\mathrm{log}(v)`$, so that $`^2=\{(y,z)^2\}`$ with $`ds^2=e^{2z}dy^2+dz^2`$.
###### Lemma 10.2.
There exist two smooth functions $`f`$ and $`g`$ on $`^2`$ such that
1. $`f`$ and $`g`$ are non negative;
2. $`f(y,z)=g(y,z)=0`$ if $`z0`$ or $`|y|1`$;
3. $`df`$ and $`dgL^r(^2,\mathrm{\Lambda }^1)`$ for any $`1<r\mathrm{}`$;
4. the support of $`dfdg`$ is contained in $`\{(y,z):|y|1,\mathrm{\hspace{0.17em}0}z1\}`$;
5. $`dfdg0`$;
6. $`{\displaystyle _^2𝑑f}dg=1`$;
7. $`{\displaystyle \frac{f}{y}}`$ and $`{\displaystyle \frac{g}{y}}L^{\mathrm{}}(^2)`$;
8. $`{\displaystyle \frac{f}{z}}`$ and $`{\displaystyle \frac{g}{z}}`$ have compact support.
Remark The forms $`df`$ and $`dg`$ cannot have compact support, otherwise, by Stokes theorem, we would have $`_^2𝑑fdg=0`$.
Proof Choose smooth functions $`h_1`$, $`h_2`$, and $`k:`$ with the following properties:
1) $`h_1,h_2`$ and $`k`$ are $`0`$;
2) $`h_i(y)=0`$ if $`|y|1`$;
3) $`h_1^{}(y)h_2(y)0`$ and $`h_1(y)h_2^{}(y)0`$ for all $`y`$;
4) the function $`(h_1^{}(y)h_2(y)h_1(y)h_2^{}(y))`$ has non empty support;
5) $`k^{}(z)0`$ for all $`z`$;
6) $`k(z)=1`$ if $`z1`$ and $`k(z)=0`$ if $`z0`$.
We set $`f(y,z):=h_1(y)k(z)`$ and $`g(y,z):=h_2(y)k(z)`$. Properties (1) and (2) of the lemma are then clear. We prove (3) (i.e. that $`dfL^r`$ for any $`1<r\mathrm{}`$).
Indeed,
$$df=h_1(y)k^{}(z)dz+k(z)h_1^{}(y)dy.$$
The first term $`h_1(y)k^{}(z)dz`$ has compact support, and the second term $`k(z)h_1^{}(y)dy`$ has its support in the infinite rectangle $`Q=\{|y|1z0\}`$.
Choose $`D<\mathrm{}`$ such that $`|k(z)h_1^{}(y)|D`$ on $`\mathrm{\Omega }`$. We have
$$|k(z)h_1^{}(y)dy|D|dy|=De^z,$$
thus, since the element of area of $`^2`$ is $`dA=e^zdydz`$, we have
$$_^2|k(z)h_1^{}(y)dy|^r𝑑AD^r_Qe^{rz}e^z𝑑y𝑑z2CD^r_0^{\mathrm{}}e^{(1r)z}𝑑z<\mathrm{},$$
from which one gets $`dfL^r`$.
Now observe that
$$dfdg=((k(z)k^{}(z))(h_1^{}(y)h_2(y)h_1(y)h_2^{}(y))dydz,$$
hence the properties (4) and (5) follow from the construction of $`h_1,h_2`$ and $`k`$.
Property (6) is only a normalization. It can be achieved by multiplying $`f`$ (or $`g`$) by a suitable constant.
Properties (7) and (8) are easy to check.
Proof of Theorem 10.1 Define the $`1`$-forms $`\alpha =df`$ and $`\gamma =dg`$ on $`^2`$ (where $`f`$ and $`g`$ are as in Lemma 10.2). It is clear that $`d\alpha =d\gamma =0`$. We also know that $`\alpha L^p`$ for any $`1<p<\mathrm{}`$ and that $`\gamma `$ is smooth and $`\gamma L^p^{}L^q^{}`$ for all $`1<p^{},q^{}<\mathrm{}`$.
Since $`{\displaystyle _^2}\alpha \gamma 0`$, we see by proposition 8.4 that $`\alpha \overline{B}_{q,p}^1(^2)`$.
Now using the isometry group of $`^2`$, we produce an infinite family of linearly independent classes in $`\overline{H}_{q,p}^1(^2)`$. ∎
## 11. The cohomology of the ball
Since the unit ball $`𝔹^n^n`$ has finite volume, we have for all $`1p,q\mathrm{}`$ $`H_{q,p}^0(𝔹^n)=\overline{H}_{q,p}^0(𝔹^n)=`$.
In higher degree, the vanishing of the De Rham cohomology of $`𝔹^n`$ is traditionally called the Poincaré Lemma; it is proved by explicitly constructing a primitive to any closed form. To prove the vanishing of the $`L_{q,p}`$cohomology of the ball, we need to control the $`L^q`$ norm of the primitive of a closed $`L^p`$-norm. For the case $`p=q`$, this was done by Gol’dshtein, Kuz’minov and Shvedov in \[8, Lemma 3.2\] and for more general $`q`$ by Iwaniec and Lutoborski in . They proved the following
###### Theorem 11.1.
For any bounded convex domain $`U^n`$ and any $`k=1,2,\mathrm{},n`$, there exists an operator
$$T=T_U:L_{loc}^1(U,\mathrm{\Lambda }^k)L_{loc}^1(U,\mathrm{\Lambda }^{k1})$$
with the following properties:
1. $`T(d\theta )+dT\theta =\theta `$ (in the sense of currents);
2. $`\left|T\theta (x)\right|C{\displaystyle _U}{\displaystyle \frac{|\theta (y)|}{|yx|^{n1}}}𝑑y`$.
###### Corollary 11.2.
The operator $`T`$ maps $`L^p(U,\mathrm{\Lambda }^k)`$ continuously to $`L^q(U,\mathrm{\Lambda }^{k1})`$ in the following cases:
either
1. $`1p,q\mathrm{}`$ and $`\frac{1}{p}\frac{1}{q}<\frac{1}{n}`$,
or
1. $`1<p,q\mathrm{}`$ and $`\frac{1}{p}\frac{1}{q}\frac{1}{n}`$.
Remark Note that condition (i) is equivalent to $`pn`$ or $`p<n`$ and $`q<\frac{np}{np}`$ and condition (ii) is relevant to conformal cohomology $`\frac{1}{p_k}\frac{1}{p_{k1}}=\frac{1}{n}`$.
Proof Assume first that $`\frac{1}{p}\frac{1}{q}<\frac{1}{n}`$ and recall the Young inequality for convolution (see \[5, Prop. 8.9\]), which says that if $`1r,s,t\mathrm{}`$ satisfy $`\frac{1}{r}+\frac{1}{s}=1+\frac{1}{t}`$, then $`fg_{L^t}f_{L^r}g_{L^s}`$. Applying this inequality to $`f=|\theta |`$ and $`g=|x|^{1n}`$ with $`r=p`$, $`t=q`$ and $`s=\frac{pq}{p+pqq}`$, and observing that
$$\frac{1}{p}\frac{1}{q}<\frac{1}{n}s(1n)>ng_{L^s(U)}<\mathrm{},$$
we conclude from previous proposition that $`T:L^p(U,\mathrm{\Lambda }^k)L^q(U,\mathrm{\Lambda }^{k1})`$ is bounded with norm at most $`|x|^{1n}_{L^s(U)}`$.
If $`p>1`$ and $`\frac{1}{p}\frac{1}{q}=\frac{1}{n}`$, then the conclusion also holds by the Hardy-Litlewood-Sobolev inequality (see \[16, p. 119\]).
###### Corollary 11.3.
The operator $`T:\mathrm{\Omega }_{p,r}^k(U)\mathrm{\Omega }_{q,p}^{k1}(U)`$ is bounded and for any $`\omega \mathrm{\Omega }_{p,r}^k(U)`$ we have $`Td\omega +dT\omega =\omega `$ provided either
i) $`1p,q,r\mathrm{}`$ such that $`\frac{1}{p}\frac{1}{q}<\frac{1}{n}`$ and $`\frac{1}{r}\frac{1}{p}<\frac{1}{n}`$,
or
ii) $`1<p,q,r\mathrm{}`$ such that $`\frac{1}{p}\frac{1}{q}\frac{1}{n}`$ and $`\frac{1}{r}\frac{1}{p}\frac{1}{n}`$.
Proof The proof is immediate from the previous Theorem and Corollary.
The Corollary 11.2 implies the following Poincaré Lemma :
###### Proposition 11.4.
Suppose that $`p,q`$ satisfy either
1. $`1p,q\mathrm{}`$ and $`\frac{1}{p}\frac{1}{q}<\frac{1}{n}`$,
or
1. $`1<p,q\mathrm{}`$ and $`\frac{1}{p}\frac{1}{q}\frac{1}{n}`$.
Then $`H_{q,p}^k(𝔹^n)=0`$ for any $`k=1,\mathrm{},n`$.
Proof Let $`\omega `$ be an arbitrary element in $`Z_p^k(𝔹^n)`$. By Corollary 11.2, we have $`T\omega L^q(𝔹^n,\mathrm{\Lambda }^{k+1})`$, since $`\omega =dT\omega +Td\omega =d(T\omega )`$ we conclude that $`[\omega ]=0H_{q,p}^k(𝔹^n)`$ and thus $`H_{q,p}^k(𝔹^n)=0`$.
If $`p,q>1`$, we have a necessary and sufficient condition :
###### Theorem 11.5.
If $`1<p,q\mathrm{}`$ and $`k=1,\mathrm{},n`$, then $`H_{q,p}^k(𝔹^n)=0`$ if and only if $`\frac{1}{p}\frac{1}{q}\frac{1}{n}`$.
Proof We know from the previous Proposition that the condition is sufficient .
To prove that $`H_{q,p}^k(𝔹^n)0`$ if $`p<n`$ and $`q>\frac{np}{np}`$, we will use Proposition 8.2. Let us fix a number $`\mu `$ in the interval $`k\frac{n}{p}<\mu <k1\frac{n}{q}`$ (which is possible since $`\frac{1}{p}>\frac{1}{q}+\frac{1}{n}`$); and choose two forms $`\theta C^{\mathrm{}}(𝕊^{n1},\mathrm{\Lambda }^{k1})`$ and $`\phi C^{\mathrm{}}(𝕊^{n1},\mathrm{\Lambda }^{nk1})`$ such that
$$_{𝕊^{n1}}\phi d\theta =1.$$
For any $`0<t<1/4`$, we choose a smooth function $`h_t:`$ such that $`h(t,r)=0`$ if $`r<t`$ or $`r>1t`$ and $`h(t,r)=\frac{1}{\left|\mathrm{log}2t\right|}`$ if $`r<12t`$ or $`r>2t`$.
Let us then consider the forms
$`\alpha `$ $`:=`$ $`d\left(r^\mu \theta \right)`$
$`\gamma _t`$ $`:=`$ $`h_t(r)r^{(\mu +1)}dr\phi `$
Step 1 The form $`\alpha `$ belongs to $`L^p(𝔹^n,\mathrm{\Lambda }^k)`$.
We will use the same notation $`\theta `$ and $`\phi `$ for a pullback of corresponding forms from $`𝕊^n`$ to $`𝔹^n\{0\}`$ induced by the radial projection in polar coordinates.
We have
$$\alpha =r^\mu \left(d\theta +\mu \frac{1}{r}dr\theta \right).$$
Because $`\left|\theta \right|r^{(k1)}`$ and $`\left|d\theta \right|r^k`$ we have $`|\alpha |r^{\mu k}`$. Therefore
$$_{𝔹^n}|\alpha |^p𝑑x_0^1\left(r^{\mu k}\right)^pr^{n1}𝑑r<\mathrm{}$$
because $`p(\mu k)+n1>p(k\frac{n}{p}k)+n1>1`$.
Step 2 The quantity $`\left|_{𝔹^n}\alpha \gamma _t\right|`$ is bounded below.
We have $`\alpha \gamma _t=h_t(r)r^1dr\phi d\theta `$; since $`_{𝕊^{n1}}\phi d\theta =1`$, we have by Fubini Theorem
$$\left|_{𝔹^n}\alpha \gamma _t\right|=_0^1h_t(r)r^1𝑑r\frac{1}{|\mathrm{log}2t|}_{2t}^{12t}r^1𝑑r1$$
as $`t0`$. This implies that $`\left|_{𝔹^n}\alpha \gamma _t\right|`$ is bounded below for small values of $`t`$.
Step 3 We have $`d\gamma _t_{L^q^{}(𝔹^n)}0`$ as $`t0`$:
We have $`d\gamma _t:=h_t(r)r^{(\mu +1)}dr\phi `$ with $`0h_t\frac{1}{|\mathrm{log}2t|}`$. Since $`\left|dr\phi \right|r^{n+k}`$, we have
$$|d\gamma _t|\frac{r^{\mu 1+kn}}{|\mathrm{log}2t|}$$
and by Fubini Theorem
$`{\displaystyle _{𝔹^n}}|d\gamma _t|^q^{}𝑑x`$ $`=`$ $`{\displaystyle _{𝔹^n}}|h_t(r)r^{(\mu +1)}dr\phi |^q^{}𝑑x`$
$``$ $`\left({\displaystyle \frac{1}{|\mathrm{log}2t|}}\right)^q^{}{\displaystyle _0^1}\left(r^{\mu 1+kn}\right)^q^{}r^{n1}𝑑r.`$
Because
$$q^{}(\mu 1+kn)+n=q^{}(\mu 1+kn(1\frac{1}{q^{}}))=q^{}(\mu 1+k\frac{n}{q})>0$$
we have
$$_0^1\left(r^{\mu 1+kn}\right)^q^{}r^{n1}𝑑r<\mathrm{}.$$
Therefore
$$\underset{t0}{lim}_{𝔹^n}|d\gamma _t|^q^{}𝑑x\underset{t0}{lim}\left(\frac{1}{|\mathrm{log}2t|}\right)^q^{}_0^1\left(r^{\mu 1+kn}\right)^q^{}r^{n1}𝑑r=0$$
Since $`\gamma _t`$ are smooth forms with compact support, Proposition 8.2 implies that $`[\alpha ]0`$ in $`H_{q,p}^k(𝔹^n)`$.
###### Corollary 11.6.
The conformal cohomology of the hyperbolic space $`^n`$ vanishes for any degree $`k>1`$, i.e.
$$H_{\frac{n}{k1},\frac{n}{k}}^k(^n)=0.$$
Proof Since the hyperbolic space $`^n`$ is conformally equivalent to the ball $`𝔹^n^n`$, this result follows at once from the conformal invariance of conformal cohomology and the previous theorem.
###### Remark 11.7.
Because $`H_{q,p}^1(^2)0`$ for any $`q,p`$, the previous corollary does not hold for $`k=1`$.
## 12. Regularization of forms and cohomology classes
In this section we investigate two different but related problems. The first one is a density result for smooth forms in $`\mathrm{\Omega }_{q,p}^{}(M)`$ and the second one is a result about representation of the cohomology $`H_{q,p}^{}(M)`$ by smooth forms. We will use the de Rham regularization method and its version for $`L_p`$-cohomology in combination with the results of section 11.
### 12.1. Regularization operators for differential forms.
The standard way of smoothing a function in $`^n`$ is by convolution with a smooth mollifier. This procedure extends to differential forms and more generally to any tensor. In his book, De Rham proposes a clever way of localizing this construction and grafting it on manifolds.
Following De Rham, we associate to any vector $`v^n`$ the map $`s_v:^n^n`$ defined by
$$s_v(x)=\{\begin{array}{cc}h^1(h(x)+v)\hfill & \text{if }x<1,\hfill \\ x\hfill & \text{if }x1.\hfill \end{array}$$
where $`h:𝔹^n^n`$ is a radial diffeomorphism such that
$$h(x)=\{\begin{array}{cc}x\hfill & \text{if }x<1/3,\hfill \\ \frac{1}{x}\mathrm{exp}(\frac{1}{(1x^2)})x\hfill & \text{if }x2/3.\hfill \end{array}$$
###### Lemma 12.1.
The map $`vs_v`$ defines an action of the group $`^n`$ on the space $`^n`$ satisfying the following properties:
1. For every $`v^n`$, the map $`s_v:^n^n`$ is a smooth diffeomorphism;
2. The mapping $`s:^n\times ^n^n`$ is smooth;
3. $`s_v`$ is the identity outside of $`𝔹^n`$;
4. For every $`x𝔹^n`$ the mapping $`v\alpha _x(v):=s_v(x)`$ is a diffeomorphism of $`^n`$ onto $`𝔹^n`$.
Proof For the first two assertions, see . The assertions (c) and (d) are obvious. ∎
Let us fix an arbitrary bounded convex domain $`U`$ such that $`\overline{𝔹}^nU^n`$. We now define the regularization operator $`R_ϵ:L_{loc}^1(U,\mathrm{\Lambda }^k)L_{loc}^1(U,\mathrm{\Lambda }^k)`$ by
$$R_\epsilon \omega :=_^ns_v^{}(\omega )\rho _\epsilon (v)𝑑v$$
where $`\rho _\epsilon (v)=\rho (v/\epsilon )`$ is a standard mollifier.
###### Proposition 12.2.
The regularization operator defined above satisfies the following properties :
1. For any $`\omega L_{loc}^1(U,\mathrm{\Lambda }^k)`$, the form $`R_ϵ\omega `$ is smooth in $`𝔹^n`$ and $`R_ϵ\omega =\omega `$ in $`U𝔹^n`$;
2. for any $`\omega \mathrm{\Omega }_{q,p}^k(U)`$, we have $`dR_\epsilon \omega =R_\epsilon d\omega `$.
3. For any $`1p,q<\mathrm{}`$ and any $`\epsilon >0`$, the operator
$$R_\epsilon :\mathrm{\Omega }_{q,p}^k(U)\mathrm{\Omega }_{q,p}^k(U)$$
is bounded and its norm satisfies $`\underset{\epsilon 0}{lim}R_\epsilon _{q,p}=1`$;
4. For any $`1p,q<\mathrm{}`$ and any $`\omega \mathrm{\Omega }_{q,p}^k(U)`$, we have
$$\underset{\epsilon 0}{lim}R_\epsilon ^{}\omega \omega _p=0.$$
Proof The first two properties are proved in . Property (3) follows from (2) and \[9, Lemma 2\] and (4) is a standard property of the regularization.
### 12.2. Homotopy operator
Given a bounded convex domain $`U^n`$ containing the closed unit ball, we introduce the homotopy
$$A_ϵ:=(IR_\epsilon )T_U:L_{loc}^1(U,\mathrm{\Lambda }^k)L_{loc}^1(U,\mathrm{\Lambda }^{k1}),$$
where $`T_U`$ is the operator defined in Theorem 11.1.
###### Lemma 12.3.
The operator $`A_\epsilon `$ is a homotopy between the Identity and the regularization operator $`R_\epsilon `$, i.e. it satisfies
$$(IR_\epsilon )\omega =dA_\epsilon \omega +A_\epsilon d\omega .$$
Proof We know from Theorem 11.1 that $`Td\omega +dT\omega =\omega `$ for all $`\omega L_{loc}^1(U,\mathrm{\Lambda }^{k1})`$, hence we have
$`dA_\epsilon \omega +A_\epsilon d\omega `$ $`=`$ $`d(IR_\epsilon )T\omega +(IR_\epsilon )Td\omega `$
$`=`$ $`dT\omega dR_\epsilon T\omega +Td\omega R_\epsilon Td\omega `$
$`=`$ $`(dT\omega +Td\omega )R_\epsilon (dT\omega +Td\omega )`$
$`=`$ $`(IR_\epsilon )(Td\omega +dT\omega )`$
$`=`$ $`(IR_\epsilon )\omega .`$
###### Proposition 12.4.
Let $`U^n`$ be a bounded convex domain containing the closed unit ball. Then $`A_\epsilon :\mathrm{\Omega }_{p,r}^k(U)\mathrm{\Omega }_{q,p}^{k1}(U)`$ is a bounded operator for any $`k=1,2,\mathrm{},n`$ in the following two cases:
i) $`1p,q,r\mathrm{}`$ such that $`\frac{1}{p}\frac{1}{q}<\frac{1}{n}`$,
ii) $`1<p,q\mathrm{}`$ and $`\frac{1}{p}\frac{1}{q}\frac{1}{n}`$ and $`\frac{1}{r}\frac{1}{p}\frac{1}{n}`$.
Furthermore, we have $`(IR_\epsilon )\omega =dA_\epsilon \omega +A_\epsilon d\omega `$ for any $`\omega \mathrm{\Omega }_{p,r}^k(U)`$ and $`A_\epsilon \omega =0`$ outside the unit ball.
Proof The first assertion follows from Proposition 12.2 and Corollary 11.3 and the second one is the previous Lemma. The last assertion follows from the fact that $`R_\epsilon =I`$ outside of the unit ball.
### 12.3. Globalization
This regularization operators $`R_\epsilon `$ and $`A_\epsilon `$ can be globalized as follow: given a Riemannian manifold $`(M,g)`$, we can find a countable atlas $`\{\phi _i:V_iMU_i\}_i`$ such that $`U_i^n`$ is a bounded convex domain satisfying $`\overline{𝔹}^nU_i^n`$ for all $`i`$ and that $`\{B_i\}`$ is a covering of $`M`$, where $`B_i:=\phi _i^1(B)V_i`$. We also assume that $`\{V_i\}`$ (and hence $`\{B_i\}`$) is a locally finite covering of $`M`$ (we can in fact assume that any collection of $`n+2`$ different charts $`V_i`$ has an empty intersection, where $`n=dimM`$.)
For any $`m`$, we define two operators
$$R_\epsilon ^{(m)},A_\epsilon ^{(m)}:L_{loc}^1(M,\mathrm{\Lambda }^m)L_{loc}^1(M,\mathrm{\Lambda }^m)$$
as follow:
$$R_\epsilon ^{(m)}:=R_{1,\epsilon }R_{2,\epsilon }\mathrm{}R_{m,\epsilon },$$
and
$$A_\epsilon ^{(m)}:=R_{1,\epsilon }R_{2,\epsilon }\mathrm{}R_{m1,\epsilon }A_{m,\epsilon },$$
where
$$R_{i,\epsilon }(\theta ):=\left(\phi _i^1\right)^{}R_\epsilon \phi _i^{}(\theta );$$
and
$$A_{i,\epsilon }(\theta ):=\left(\phi _i^1\right)^{}(R_{i,\epsilon }I)T_{U_i}\phi _i^{}(\theta ).$$
Here $`T_{U_i}`$ is the operator defined on the domain $`U_i`$ in Theorem 11.1.
Observe that the operator $`R_{i,\epsilon }`$ is a priori only defined on $`V_i`$, but it acts as the identity on $`V_i\overline{B}_i`$ and can thus be extended on the whole of $`M`$ by declaring that $`R_{i,\epsilon }=id`$ on $`M\overline{B}_i`$. Likewise, the operator $`A_{i,\epsilon }`$ is a priori only defined on $`V_i`$, but it is zero on $`V_i\overline{B}_i`$ (because $`R_\epsilon =I`$ outside of the unit ball). Hence $`A_{i,\epsilon }`$ can be extended on the whole of $`M`$ by declaring $`A_{i,\epsilon }=0`$ on $`M\overline{B}_i`$.
We now define the global regularization operator and the global homotopy operator as follow:
(12.1)
$$R_\epsilon ^M:=\underset{m\mathrm{}}{lim}R_\epsilon ^{(m)},A_\epsilon ^M:=\underset{m=1}{\overset{\mathrm{}}{}}A_\epsilon ^{(m)}.$$
By construction, the expressions $`R_\epsilon ^M:=_iR_{i,\epsilon }`$ and $`A_\epsilon ^M:=_lA_\epsilon ^{(k)}`$ are really finite operations in any compact set and the operators $`R_\epsilon ^M,A_\epsilon ^M`$ are thus well defined on $`L_{loc}^1(M,\mathrm{\Lambda }^k)`$.
###### Theorem 12.5.
For every Riemannian manifold $`M`$ there exists a family of regularization operators $`R_\epsilon ^M`$ and homotopy operators $`A_\epsilon ^M`$ such that
1. For any $`\omega L_{loc}^1(M,\mathrm{\Lambda }^k)`$, the form $`R_ϵ^M\omega `$ is smooth in $`M`$;
2. For any $`\omega \mathrm{\Omega }_{q,p}^k(M)`$, we have $`dR_\epsilon ^M\omega =R_\epsilon ^Md\omega `$;
3. For any $`1p,q<\mathrm{}`$ and any $`\epsilon >0`$, the operator $`R_\epsilon ^M:\mathrm{\Omega }_{q,p}^k(M)\mathrm{\Omega }_{q,p}^k(M)`$ is bounded and its norm satisfies $`\underset{\epsilon 0}{lim}R_\epsilon ^M_{q,p}=1`$;
4. For any $`1p,q<\mathrm{}`$ and any $`\omega \mathrm{\Omega }_{q,p}^k(M)`$ we have
$$\underset{\epsilon 0}{lim}R_\epsilon ^M\omega \omega _p=0.$$
5. The operator $`A_\epsilon :\mathrm{\Omega }_{pr}^k(M)\mathrm{\Omega }_{q,p}^{k1}(M)`$ is bounded for any $`k=1,\mathrm{},n`$ in the following cases:
1. $`1p,q,r\mathrm{}`$ such that $`\frac{1}{p}\frac{1}{q}<\frac{1}{n}`$ and $`\frac{1}{r}\frac{1}{p}<\frac{1}{n}`$,
2. $`1<p,q,r\mathrm{}`$ such that $`\frac{1}{p}\frac{1}{q}\frac{1}{n}`$ and $`\frac{1}{r}\frac{1}{p}\frac{1}{n}`$.
6. We have the homotopy formula
$$\omega R_\epsilon ^M\omega =dA_\epsilon ^M\omega +A_\epsilon ^Md\omega .$$
Proof The first four assertions follow immediately from Proposition 12.2.
The fifth assertion follows from Proposition 12.2 and Corollary 11.3.
To prove the last assertion, observe that by Lemma 12.3, we have $`\omega R_{m,\epsilon }\omega =dA_{m,\epsilon }\omega +A_{m,\epsilon }d\omega `$. Multiplying this expression by $`R_\epsilon ^{(m1)}`$, we obtain
$$R_\epsilon ^{(m1)}\omega R_\epsilon ^{(k)}\omega =dA_\epsilon ^{(k)}\omega +A_\epsilon ^{(m)}d\omega ,$$
summing this identities on $`m=1,2,\mathrm{}`$, we obtain the assertion (6).
###### Corollary 12.6.
For any $`q,p[1,\mathrm{})`$, the space
$$C^{\mathrm{}}\mathrm{\Omega }_{q,p}^k(M):=C^{\mathrm{}}(M)\mathrm{\Omega }_{q,p}^k(M)$$
of smooth $`k`$-forms $`\theta `$ in $`L^p`$ such that $`d\theta L^q`$ is dense in $`\mathrm{\Omega }_{q,p}^k(M)`$.
Proof This result follows immediately from the first three conditions in Theorem 12.5. ∎
### 12.4. $`L_\pi `$-cohomology and smooth forms
The previous theorem implies that under suitable assumptions on $`p,q`$, the $`L_\pi `$-cohomology of a Riemannian manifold can be represented by smooth forms.
To be more precise, for any sequence $`\pi `$, we denote by
$$C^{\mathrm{}}\mathrm{\Omega }_\pi ^k(M):=C^{\mathrm{}}(M)\mathrm{\Omega }_\pi ^k(M)$$
the subcomplex of smooth forms in $`\mathrm{\Omega }_\pi ^k(M)`$ and by
$$C^{\mathrm{}}H_\pi ^{}(M)=H^{}(C^{\mathrm{}}\mathrm{\Omega }_\pi ^k(M))$$
its cohomology.
###### Theorem 12.7.
Let $`(M,g)`$ be a $`n`$-dimensional Riemannian manifold and
$`\pi =\{p_0,p_1,\mathrm{},p_n\}(1,\mathrm{})`$ a finite sequence of numbers such that
$`\frac{1}{p_k}\frac{1}{p_{nk}}\frac{1}{n}`$ for $`k=1,2,..n`$. Then
$$C^{\mathrm{}}H_\pi ^{}(M)=H_\pi ^{}(M).$$
Proof This result follows immediately from Proposition 4.8 and Theorem 12.5.
It is perhaps useful to reformulate this theorem without the language of complexes:
###### Theorem 12.8.
Let $`(M,g)`$ be a $`n`$-dimensional Riemannian manifold and suppose that $`p,q(1,\mathrm{})`$ satisfy $`\frac{1}{p}\frac{1}{q}\frac{1}{n}`$. Then the cohomology $`H_{q,p}^{}(M)`$ can be represented by smooth forms.
More precisely, any closed form in $`Z_p^k(M)`$ is cohomologous to a smooth form in $`L^p(M)`$. Furthermore, if two smooth closed forms $`\alpha ,\beta C^{\mathrm{}}(M)Z_p^k(M)`$ are cohomologous modulo $`d\mathrm{\Omega }_{q,p}^{k1}(M)`$, then they are cohomologous modulo $`dC^{\mathrm{}}\mathrm{\Omega }_{q,p}^{k1}(M)`$.
###### Corollary 12.9.
Let $`(M,g)`$ be a $`n`$-dimensional Riemannian manifold and suppose that $`p,q(1,\mathrm{})`$ satisfy $`\frac{1}{p}\frac{1}{q}\frac{1}{n}`$. Then any reduced cohomology class can be represented by a smooth form.
Proof This is clear from the previous Theorem, since $`\overline{H}_{q,p}^k(M)`$ is a quotient of $`H_{q,p}^k(M)`$.
### 12.5. The case of compact manifolds
From previous results, we now immediately have:
###### Theorem 12.10.
Let $`(M,g)`$ be a compact $`n`$-dimensional Riemannian manifold and $`\pi =\{p_0,p_1,\mathrm{},p_n\}(1,\mathrm{})`$ a finite sequence of numbers such that $`\frac{1}{p_k}\frac{1}{p_{nk}}\frac{1}{n}`$ for $`k=1,2,..n`$. Then
$$H_\pi ^{}(M)=H_{\mathrm{DeRham}}^{}(M).$$
In particular $`H_\pi ^{}(M)`$ is finite dimensional and thus $`T_\pi ^{}(M)=0`$.
Proof Recall that the De Rham cohomology $`H_{\mathrm{DeRham}}^{}(M)`$ of $`M`$ is the cohomology of the complex $`(C^{\mathrm{}}(M,\mathrm{\Lambda }^{}),d)`$. Any smooth form on a compact Riemannian manifold clearly belongs to $`L^p`$ for any $`p[0,\mathrm{}]`$, hence $`(C^{\mathrm{}}(M,\mathrm{\Lambda }^{}),d)=C^{\mathrm{}}\mathrm{\Omega }_\pi ^k(M)`$ and by Theorem 12.7, we have
$$H_\pi ^{}(M)=C^{\mathrm{}}H_\pi ^{}(M)=H_{\mathrm{DeRham}}^{}(M).$$
It is well known that the De Rham cohomology of a compact manifold is finite dimensional. Since $`dimT_\pi ^{}(M)dimH_\pi ^{}(M)<\mathrm{}`$, it follows from Lemma 4.4 that $`T_\pi ^{}(M)=0`$.
### 12.6. Proof of Theorems 1.1 and 1.2
Let us define the sequence
$`\pi =\{p_0,p_1,\mathrm{},p_n\}`$ by $`p_j=q`$ if $`j=1,2,..k1`$ and $`p_j=p`$ if $`j=k,\mathrm{},n`$.
By hypothesis, we have $`\frac{1}{p}\frac{1}{q}\frac{1}{n}`$, hence the sequence $`\pi `$ satisfies $`\frac{1}{p_j}\frac{1}{p_{j1}}\frac{1}{n}`$ for all $`j`$. Hence we know by Theorem 12.10 that $`H_{q,p}^k(M)=H_{\mathrm{DeRham}}^k(M)`$ and $`T_{q,p}^k(M)=0`$.
Thus Theorem 1.1 follows from Theorem 6.2 and Theorem 1.2 follows from Theorem 6.1.
## 13. Relation with a non linear PDE
We show in this section that the vanishing of torsion gives sufficient condition to solving the non linear equation
(13.1)
$$\delta (d\theta ^{p2}d\theta )=\alpha ,$$
where $`\delta `$ is the operator defined for $`\omega L_{loc}^1(M,\mathrm{\Lambda }^k)`$ as
$$\delta \omega =(1)^{nk+n+1}d\omega .$$
Recall that for any $`k`$-form $`\omega `$, we have<sup>2</sup><sup>2</sup>2Here is the proof: Since $`\omega `$ is a $`k`$ form, $`d\omega `$ is a form of degree $`m=nk+1`$ and $`d\omega =(1)^{m(nm)}d\omega =(1)^{nk+n+1+k}d\omega `$, therefore $`(1)^kd\omega =(1)^{nk+n+1}d\omega =\delta \omega `$..
(13.2)
$$\delta \omega =(1)^kd\omega .$$
This operator is the formal adjoint to the exterior differential $`d`$ in the sense that
(13.3)
$$_M\omega ,d\phi 𝑑\mathrm{vol}=_M\delta \omega ,\phi 𝑑\mathrm{vol}$$
for any $`\phi C_c^{\mathrm{}}(M,\mathrm{\Lambda }^{k1})`$.
Indeed, by definition of the Hodge $``$ operator, we have
$$d\phi ,\omega d\mathrm{vol}=(d\phi \omega )$$
and from the definition of the weak exterior differential, it follows that
$$_Md\phi ,\omega d\mathrm{vol}=_Md\phi \omega =(1)^k_M\phi d\omega .$$
Thus from (13.2):
$`{\displaystyle _M}d\phi ,\omega 𝑑\mathrm{vol}`$ $`=`$ $`(1)^k{\displaystyle _M}\phi d\omega `$
$`=`$ $`{\displaystyle _M}\phi \delta \omega `$
$`=`$ $`{\displaystyle _M}\phi ,\delta \omega 𝑑\mathrm{vol}.`$
Applying (13.3) to $`\omega =|d\theta |^{p2}d\theta `$, we obtain the following
###### Lemma 13.1.
$`\theta L_{loc}^1(M,\mathrm{\Lambda }^k)`$ is a solution to (13.1) if and only if
(13.4)
$$_Md\phi ,d\theta ^{p2}d\theta 𝑑\mathrm{vol}=_M\phi ,\alpha 𝑑\mathrm{vol}$$
for any $`\phi C_c^{\mathrm{}}(M,\mathrm{\Lambda }^k)`$ .
The equation (13.4) is just the weak form of (13.1).
Remark In the scalar case, equation (13.1) is just the $`p`$-Laplacian. The case of differential forms on the manifold $`M=^n`$ appears in section 6.1 of where it is investigated by the method of Hodge dual systems, see also \[12, §8\].
###### Theorem 13.2.
Assume $`T_{q,p}^k(M)=0`$, $`(1<q,p<\mathrm{})`$ and $`\alpha L^q^{}(M,\mathrm{\Lambda }^k)`$ where $`q^{}=q/(q1)`$.
(A) If $`{\displaystyle _M}\alpha ,\phi 𝑑\mathrm{vol}=0`$ for any $`\phi Z_q^k(M)`$, then (13.4) has a solution $`\theta \mathrm{\Omega }_{q,p}^k(M)`$.
(B) Conversely, if (13.4) is solvable in $`\mathrm{\Omega }_{q,p}^k(M)`$, then $`{\displaystyle _M}\alpha ,\phi 𝑑\mathrm{vol}=0`$ for any $`\phi C_c^{\mathrm{}}(M,\mathrm{\Lambda }^k)`$ such that $`d\phi =0`$.
Proof Assertion (B) follows from the previous Lemma, because for any $`\phi C_c^{\mathrm{}}(M,\mathrm{\Lambda }^k)\mathrm{ker}d`$, we have
$$_M\alpha ,\phi 𝑑\mathrm{vol}=_Md\theta ^{p2}d\theta ,d\phi 𝑑\mathrm{vol}=0.$$
Let us prove assertion (A). The variational functional corresponding to (13.4) reads
$$I(\theta )=\frac{1}{p}_Md\theta ^p𝑑\mathrm{vol}_M\alpha ,\theta 𝑑\mathrm{vol}.$$
We first show that the functional $`I(\theta ):\mathrm{\Omega }_{q,p}^k(M)`$ is bounded from below:
For any $`\theta \mathrm{\Omega }_{q,p}^k(M)`$ there exists a unique element $`z_q(\theta )Z_q^k(M)`$ such that $`\theta z_q(\theta )_qinf_{zZ_q^k(M)}\theta z_q`$; this follows from the uniform convexity of $`\mathrm{\Omega }_{q,p}^k(M)`$. Since $`T_{q,p}^k(M)=0`$, the Proposition 1.2 implies that
(13.5)
$$\theta z_q(\theta )_qCd\theta _p$$
for some positive constant $`C`$. Using this inequality and Hölder’s inequality, we obtain
$$I(\theta )\frac{1}{p}d\theta _p^p\alpha _q^{}\theta z_q(\theta )_q\frac{1}{p}d\theta _p^pC\alpha _q^{}d\theta _p.$$
Since the function $`f:`$ defined by $`f(x)=\frac{1}{p}|x|^pax`$ is bounded below for $`x0`$, the previous inequality implies that
$$\underset{\theta \mathrm{\Omega }_{q,p}^k(M)}{inf}I(\theta )>\mathrm{}.$$
We now prove the existence of a minimizer of $`I`$ on $`\mathrm{\Omega }_{q,p}^k(M)`$: Let $`\left\{\theta _i\right\}\mathrm{\Omega }_{q,p}^k(M)`$ be a sequence such that $`I(\theta _i)infI(\theta )`$. Because the function $`f(x)=\frac{1}{p}|x|^pax`$ is proper, the inequality
$$I(\theta _i)\frac{1}{p}d\theta _i_p^pC\alpha _q^{}d\theta _i_p$$
implies that $`\{d\theta _i_p\}`$ is bounded and, by (13.5), $`\{\theta _iz_q(\theta _i)_q\}`$ is also bounded. Hence the sequence $`\{\stackrel{~}{\theta }_i:=\theta _iz_q(\theta _i)\}`$ is bounded in $`\mathrm{\Omega }_{q,p}^k(M)`$.
Since $`\mathrm{\Omega }_{q,p}^k(M)`$ is reflexive there exists a subsequence (still noted $`\{\stackrel{~}{\theta }_i\}`$) which converges weakly to some $`\theta _0\mathrm{\Omega }_{q,p}^k(M)`$. By the weak continuity of the functional $`_M\alpha ,\theta 𝑑\mathrm{vol}`$ in $`\mathrm{\Omega }_{q,p}^k(M)`$ we have
(13.6)
$$\underset{i\mathrm{}}{lim}_M\alpha ,\stackrel{~}{\theta }_i𝑑\mathrm{vol}=_M\alpha ,\theta _0𝑑\mathrm{vol}$$
The lower semicontinuity of the norm under the weak convergence implies that
$$d\theta _0_p\underset{i\mathrm{}}{lim\; inf}d\stackrel{~}{\theta }_i_p.$$
Combining the last inequality with (13.6) we obtain
$$I(\theta _0)\underset{i\mathrm{}}{lim\; inf}I(\theta _i)$$
and by the choice of $`\theta _i`$ we finally have $`I(\theta _0)=infI(\theta )`$.
It is now clear that $`\theta _0`$ is a solution of (13.4), hence a weak solution of (13.1).
Definition. The Riemannian manifold $`(M,g)`$ is *$`s`$-parabolic* if for any $`\epsilon >0`$, there exists a smooth function $`f_\epsilon `$ with compact support, such that $`f_\epsilon =1`$ on the ball $`B(x_0,1/\epsilon )`$ and $`df_\epsilon _{L^s(M)}\epsilon `$. where $`x_0M`$ is a fixed base point.
Some basic facts about this notion can be found in .
###### Corollary 13.3.
Assume as above that $`T_{q,p}^k(M)=0`$ and $`\alpha L^q^{}(M,\mathrm{\Lambda }^k)`$ where $`q^{}=q/(q1)`$ , $`(1<q,p<\mathrm{})`$.
Assume furthermore that $`M`$ is $`s`$-parabolic for $`\frac{1}{s}=\frac{1}{p}+\frac{1}{q}`$.
Then equation (13.4) is solvable in $`\mathrm{\Omega }_{q,p}^k(M)`$, if and only if $`_M\alpha ,\phi 𝑑\mathrm{vol}=0`$ for any $`\phi Z_q^k(M)`$.
Proof The condition is sufficient by the previous theorem. Now let $`\phi Z_q^k(M)`$ be arbitrary and let $`R_\epsilon ^M`$ be the smoothing operator and $`f_\epsilon `$ be as in the previous definition. Then
$$\phi _\epsilon :=f_\epsilon R_\epsilon ^M(\phi )C_c^{\mathrm{}}(M,\mathrm{\Lambda }^k).$$
Let us observe that
$$|d\theta |^{p2}d\theta _{L^p^{}(M)}=d\theta _{L^p(M)}^{p/p^{}}$$
where $`p^{}=p/(p1)`$. Since $`\frac{1}{s}=1\frac{1}{p^{}}+\frac{1}{q}`$, we have by Hölder’s inequality:
$`{\displaystyle _M}\alpha ,\phi _\epsilon 𝑑\mathrm{vol}`$ $`=`$ $`{\displaystyle _M}d\theta ^{p2}d\theta ,d\phi _\epsilon 𝑑\mathrm{vol}`$
$`=`$ $`{\displaystyle _M}d\theta ^{p2}d\theta ,df_\epsilon R_\epsilon ^M(\phi )𝑑\mathrm{vol}`$
$``$ $`|d\theta |^{p2}d\theta _{L^p^{}(M)}df_\epsilon _{L^s(M)}R_\epsilon ^M(\phi )_{L^q(M)}`$
$``$ $`\left(d\theta _{L^p(M)}^{p^{}/p}R_\epsilon ^M(\phi )_{L^q(M)}\right)df_\epsilon _{L^s(M)}`$
As $`\epsilon 0`$, we have $`df_\epsilon _{L^s(M)}0`$ while $`\left(d\theta _{L^p(M)}^{p^{}/p}R_\epsilon ^M(\phi )_{L^q(M)}\right)`$ remains bounded. On the other hand,
$$\underset{\epsilon 0}{lim}_M\alpha ,\phi _\epsilon 𝑑\mathrm{vol}=_M\alpha ,\phi 𝑑\mathrm{vol}$$
and the result follows. ∎
## 14. Torsion in $`L_2`$-cohomology and the Hodge-Kodaira decomposition
In this section, we study some connection between the torsion in $`L_2`$-cohomology and the Laplacian $`\mathrm{\Delta }`$ acting on differential forms on the complete Riemannian manifold $`(M,g)`$.
Recall that $`\mathrm{\Delta }=d\delta +\delta d`$ where $`\delta `$ is the formal adjoint operator to the exterior differential $`d`$. We look at $`\mathrm{\Delta }`$ as an unbounded operator acting on the Hilbert space $`L^2(M,\mathrm{\Lambda }^k)`$. In particular, all function spaces appearing in this section are subspaces of $`L^2(M,\mathrm{\Lambda }^k)`$. We denote by $`_2^k(M)=L^2(M,\mathrm{\Lambda }^k)\mathrm{ker}\mathrm{\Delta }`$ the space of $`L^2`$ harmonic forms.
We begin with the following result, which can be proved by standard arguments from functional analysis:
###### Theorem 14.1.
For any complete Riemannian manifold $`(M,g)`$, the following conditions are equivalent:
1. $`\mathrm{Im}\mathrm{\Delta }`$ is a closed subspace in $`L^2(M,\mathrm{\Lambda }^k)`$;
2. $`\mathrm{Im}\mathrm{\Delta }=\left(_2^k(M)\right)^{}`$;
3. There exists a bounded linear operator $`G:L^2(M,\mathrm{\Lambda }^k)L^2(M,\mathrm{\Lambda }^k)`$ such that for any $`\alpha L^2(M,\mathrm{\Lambda }^k)`$ we have
$$\mathrm{\Delta }G\alpha =G\mathrm{\Delta }\alpha =\alpha H\alpha $$
where $`H:L^2(M,\mathrm{\Lambda }^k)_2^k(M)`$ is the orthogonal projection onto the space of $`L^2`$ harmonic forms.
Remark: $`G`$ is called the *Green operator*. It is not difficult to check that $`dG=Gd`$ and $`\delta G=G\delta `$.
For the convenience of the reader, we briefly explain the proof of this Theorem:
Proof (a) $``$ (b): Because $`\mathrm{\Delta }`$ is self-adjoint, we know by standard functional analysis (see e.g. , page 28) that $`\overline{\mathrm{Im}\mathrm{\Delta }}=\left(_2^k(M)\right)^{}`$,
(b) $``$ (c): This follows from the Banach Open Mapping Theorem. More precisely, let us denote by
$$E:=\{\omega L^2(M,\mathrm{\Lambda }^k)|\omega _2^k(M)\mathrm{and}\mathrm{\Delta }\omega L^2(M,\mathrm{\Lambda }^k)\}$$
the domain of the Laplacian. This is a Hilbert space for the graph norm $`\omega _E:=\omega _{L^2}+\mathrm{\Delta }\omega _{L^2}`$ and the map $`\mathrm{\Delta }:E\mathrm{Im}\mathrm{\Delta }=\left(_2^k(M)\right)^{}`$ is a continuous bijective operator.
From the Banach Open Mapping Theorem, we know that the map
$$G:=\mathrm{\Delta }^1(1H):L^2(M,\mathrm{\Lambda }^k)L^2(M,\mathrm{\Lambda }^k)$$
given by the composition
$$L^2(M,\mathrm{\Lambda }^k)\stackrel{1H}{}\left(_2^k(M)\right)^{}\stackrel{\mathrm{\Delta }^1}{}EL^2(M,\mathrm{\Lambda }^k)$$
is continuous. It is clear that $`G`$ satisfies the required properties.
(c) $``$ (b): Condition (c) obviously implies that $`\mathrm{Im}\mathrm{\Delta }\left(_2^k(M)\right)^{}`$. The other inclusion $`\mathrm{Im}\mathrm{\Delta }\left(_2^k(M)\right)^{}`$ always holds since $`\mathrm{\Delta }`$ is self-adjoint.
In the case of complete Riemannian manifolds, we have the following :
###### Theorem 14.2.
For any complete Riemannian manifold $`(M,g)`$, we have
$$_2^k(M)=\mathrm{ker}d\mathrm{ker}\delta L^2(M,\mathrm{\Lambda }^k),$$
and the orthogonal decomposition
$$L^2(M,\mathrm{\Lambda }^k)=\overline{\mathrm{Im}d}\overline{\mathrm{Im}\delta }_2^k(M).$$
The first part is due to Andreotti and Vesentini, the second part is the well known Hodge-Kodaira decomposition. A proof is given in \[3, Theorem 24 and 26\].
Using both previous Theorems, we can now prove the following result:
###### Theorem 14.3.
For any complete Riemannian manifold $`(M,g)`$, the following conditions are equivalent:
1. $`\mathrm{Im}\mathrm{\Delta }=\left(_2^k(M)\right)^{}`$;
2. we have the orthogonal decomposition
$$L^2(M,\mathrm{\Lambda }^k)=\mathrm{Im}d\mathrm{Im}\delta _2^k(M);$$
3. $`\mathrm{Im}d`$ and $`\mathrm{Im}\delta `$ are closed in $`L^2(M,\mathrm{\Lambda }^k)`$;
4. $`T_2^k(M)=0`$ and $`T_2^{nk}(M)=0`$.
We will also need the following
###### Lemma 14.4.
If $`T_2^k(M)=0`$, then
$$\mathrm{Im}(\delta d)=\mathrm{Im}(\delta )$$
as subsets of $`L^2(M,\mathrm{\Lambda }^k)`$.
Proof It is clear that $`\mathrm{Im}(\delta d)\mathrm{Im}(\delta )`$. To prove the other inclusion, consider an arbitrary element $`\alpha \mathrm{Im}\delta `$. Because $`\mathrm{Im}\delta \mathrm{ker}d=Z_2^k(M)`$, we know by Theorem 13.2 that we can find a form $`\theta L^2(M,\mathrm{\Lambda }^k)`$ such that $`\delta d\theta =\alpha `$. In particular $`\alpha \mathrm{Im}\delta d`$.
Remark. Using the formula $`\delta =\pm d`$, we see that this lemma also says that $`\mathrm{Im}(d\delta )=\mathrm{Im}(d)`$, provided $`T_{2,2}^{nk}(M)=0`$.
Proof of Theorem 14.3.
(i) $``$ (ii): Condition (i) is equivalent to (c) of Theorem 14.1. Hence, assuming (i), we know that any $`\alpha L^2(M,\mathrm{\Lambda }^k)`$ can be written as
$$\alpha H\alpha =\mathrm{\Delta }G\alpha =d(\delta G\alpha )+\delta (dG\alpha )$$
and the decomposition (ii) follows.
(ii) $``$ (iii): is clear from Theorem 14.2.
(iii) $``$ (vi): Follows from the definition of torsion and the formula $`\delta =\pm d`$.
(vi) $``$ (i): We know from the previous lemma and the orthogonality of $`\mathrm{Im}d`$ and $`\mathrm{Im}\delta `$ that
$$\mathrm{Im}\mathrm{\Delta }=\mathrm{Im}(d\delta +\delta d)=\mathrm{Im}(d\delta )+\mathrm{Im}(\delta d)=\mathrm{Im}(d)+\mathrm{Im}(\delta ),$$
provided $`T_2^k(M)=T_2^{nk}(M)=0`$. In particular, $`\mathrm{Im}\mathrm{\Delta }`$ is closed, since $`\mathrm{Im}d`$ and $`\mathrm{Im}\delta `$ are closed, and we conclude by Theorem 14.1 that $`\mathrm{Im}\mathrm{\Delta }=\left(_2^k(M)\right)^{}`$.
###### Corollary 14.5.
If $`(M,g)`$ is complete, then the equation $`\mathrm{\Delta }\omega =\alpha L^2(M,\mathrm{\Lambda }^k)`$ is solvable in $`L^2(M,\mathrm{\Lambda }^k)`$ for any $`\alpha _2^k(M)`$, if and only if
$$T_2^k(M)=0\mathrm{and}T_2^{nk}(M)=0.$$
The proof is immediate.
In conclusion, we formulate the following version of Hodge Theorem and Poincaré duality for $`L^2`$-cohomology:
###### Corollary 14.6.
If $`(M,g)`$ is a complete Riemannian manifold such that
$`T_2^k(M)=T_2^{nk}(M)=0`$, then
$$\overline{H}_2^k(M)=H_2^k(M)_2^k(M)_2^{nk}(M)H_2^{nk}(M)=\overline{H}_2^{nk}(M).$$
Proof The equality $`\overline{H}_2^k(M)=H_2^k(M)`$ is equivalent to $`T_2^k(M)=0`$.
From Theorem 14.3, we know that if the torsion vanishes, then
$$\mathrm{ker}d=(\mathrm{Im}\delta )^{}=\mathrm{Im}d_2^k(M),$$
i.e. $`H_2^k(M)_2^k(M)`$ by definition of cohomology.
The isomorphism $`_2^k(M)_2^{nk}(M)`$ is given by the Hodge $``$ operator and the proof now ends as it begins.
## Appendix: A “classic” proof of Theorem 1.1 in the compact case.
In this appendix, we shortly give another proof of Theorem 1.1 for compact manifolds which is based on the Hodge De-Rham theory and the regularity theory for elliptic systems, together with some techniques from functional analysis. All these tools were available 40 years ago, however, we did not find a written proof in the literature.
We start with the fact that the space of harmonic currents on a compact Riemannian manifold $`(M,g)`$ is finite dimensional and that we can construct two linear operators acting on currents on $`M`$
$$G,H:𝒟^{}(M)𝒟^{}(M),$$
and such that
1. $`\mathrm{ker}\mathrm{\Delta }=\mathrm{Im}H=\mathrm{ker}(IH)`$;
2. $`\mathrm{ker}\mathrm{\Delta }\mathrm{Im}(IH)=\{0\}`$;
3. $`\mathrm{\Delta }G=(IH)`$;
4. $`\mathrm{\Delta }(IH)=\mathrm{\Delta }`$;
5. $`dG=Gd`$.
This result is theorem 23 in , the operator $`H`$ is the projection onto the space of harmonic forms and $`G`$ is the Green operator.
Using elliptic regularity, we can prove the following theorem:
###### Theorem 14.7.
The Green operator defines a bounded linear operator
$$G:W^{m,p}(M,\mathrm{\Lambda }^k)W^{m+2,p}(M,\mathrm{\Lambda }^k)$$
for any $`m`$. Here $`W^{m,p}(M,\mathrm{\Lambda }^k)`$ is the Sobolev space of differential forms of degree $`k`$ on $`M`$ with coefficients in $`W^{m,p}`$.
Assuming this result for the time being, let us conclude the proof of Theorem 1.1. We first state the following corollary:
###### Corollary 14.8.
For any compact Riemannian manifold $`(M,g)`$, there exists a constant $`C_1`$ such that
(14.1)
$$\theta \zeta _{W^{1,p}(M)}C_1d\theta _{L^p(M)},$$
where $`\zeta :=H\theta +d\delta G\theta `$.
Proof From previous theorem, we see that $`\delta G:L^p(M,\mathrm{\Lambda }^k)W^{1,p}(M,\mathrm{\Lambda }^{k+1})`$ is a bounded operator.
Since $`\mathrm{\Delta }G=(d\delta +\delta d)G=(IH)`$, we have $`\theta \zeta =\delta dG\theta =\delta Gd\theta `$ and thus
$$\theta \zeta _{W^{1,p}(M)}=\delta Gd\theta _{W^{1,p}(M)}C_1d\theta _{L^p(M)},$$
where $`C_1`$ is the operator norm $`C_1:=\delta G_{L^pW^{1,p}}.`$
### Proof of Theorem 1.1.
The classical Sobolev embedding theorem on compact manifolds, states in particular that there is a constant $`C_2`$ such that
(14.2)
$$\omega _{L^q(M)}C_2\omega _{W^{1,p}(M)},$$
provided that conditions (1.2), are satisfied.
Combining (14.1) and (14.2) and observing that, by the Sobolev embedding theorem and (1.2), we have $`\zeta =H\theta +d\delta G\theta Z_q^k(M)`$, we obtain (1.1) with $`C=C_1C_2`$.
### Proof of Theorem 14.7
The proof is in several steps.
Step 1. The elliptic estimate for the Laplacian acting on forms on a compact manifold says that there exists a constant $`A_m`$ such that for any form $`\theta W^{m+2,p}(M,\mathrm{\Lambda }^k)`$ we have
(14.3)
$$\theta _{W^{m+2,p}(M)}A_m\left(\mathrm{\Delta }\theta _{W^{m,p}(M)}+\theta _{W^{m,p}(M)}\right).$$
This result is deep. The case $`p=2`$ is proved in proved in \[18, §6.29\], the scalar case for any $`p(0,\mathrm{})`$ can be found in \[7, §9.5\] and the general case in \[1, Chapter IV\].
Step 2. A first consequence of this estimates is the hypoellipticity of the Laplacian, i.e. the fact if $`\mathrm{\Delta }\theta `$ is a smooth form, then $`\theta `$ itself is smooth (the proof follows from a bootstrap argument based on (14.3) and the fact that $`_{m1}W^{m,p}(M)=C^{\mathrm{}}(M)`$.) It follows in particular that the Green operator $`G`$ maps smooth forms to smooth forms.
Step 3. Using (14.3), we show that for any sequence $`\{\theta _i\}W^{m+2,p}`$, we have
(14.4)
$$\mathrm{\Delta }\theta _i_{W^{m,p}(M)}\mathrm{bounded}(IH)\theta _i_{W^{m,p}(M)}\mathrm{bounded}.$$
Indeed, otherwise there exists a sequence such $`\mathrm{\Delta }\theta _i_{W^{m,p}(M)}`$ is bounded and $`(IH)\theta _i_{W^{m,p}(M)}\mathrm{}`$. Let us set
$$\phi _i:=\frac{(IH)\theta _i}{(IH)\theta _i_{W^{m,p}(M)}}W^{m+2,p}(M),$$
we then have $`\phi _i_{W^{m,p}(M)}=1`$ and
$$\underset{i\mathrm{}}{lim}\mathrm{\Delta }\phi _i_{W^{m,p}(M)}=\frac{\mathrm{\Delta }\theta _i_{W^{m,p}(M)}}{(IH)\theta _i_{W^{m,p}(M)}}=0.$$
The elliptic estimate (14.3) gives us
$$\phi _i_{W^{m+2,p}(M)}A_m\left(\mathrm{\Delta }\phi _i_{W^{m,p}(M)}+\phi _i_{W^{m,p}(M)}\right)$$
and thus $`\{\phi _i\}`$ is bounded in $`W^{m+2,p}(M)`$.
Because $`W^{m+2,p}(M)`$ is reflexive, there exists a subsequence which converges weakly in $`W^{m+2,p}(M)`$. We still denote this subsequence by $`\{\phi _i\}`$. Let $`\phi W^{m+2,p}(M)`$ be the weak limit of this subsequence, we then have by the lower semi-continuity of the norm
$$\mathrm{\Delta }\phi _{W^{m,p}(M)}\underset{i\mathrm{}}{lim\; inf}\mathrm{\Delta }\phi _i_{W^{m,p}(M)}=0,$$
hence $`\phi \mathrm{ker}\mathrm{\Delta }`$. Since we also have $`\phi \mathrm{Im}(IH)`$ we must have $`\phi =0`$.
By the compactness of the embedding $`W^{m+2,p}(M)W^{m,p}(M)`$, we may assume that this subsequence converges strongly in $`W^{m,p}(M)`$. In particular we have
$$1=\underset{i\mathrm{}}{lim}\phi _i_{W^{m,p}(M)}=\underset{i\mathrm{}}{lim}\phi _i_{W^{m,p}(M)}=0,$$
This contradiction proves (14.4).
Step 4. We now show that:
$$\mathrm{\Delta }\left(W^{m+2,p}(M)\right)\text{is closed in}W^{m,p}(M)$$
Indeed, for any $`\omega W^{m,p}(M)`$ in the closure of $`\mathrm{\Delta }\left(W^{m+2,p}\right)`$, there exists a sequence $`\{\theta _i\}W^{m+2,p}`$, such that $`\mathrm{\Delta }\theta _i\omega `$. By step 3, $`\{(IH)\theta _i\}`$ is bounded in $`W^{m,p}`$, and by (14.3), this sequence is also bounded in $`W^{m+2,p}`$ (recall that $`\mathrm{\Delta }(IH)\theta _i=\mathrm{\Delta }\theta _i`$).
By the compactness of the embedding $`W^{m+2,p}(M)W^{m,p}(M)`$, there exists a subsequence such that $`\{(IH)\theta _i\}`$ converges strongly in $`W^{m,p}`$, and by (14.3) again, $`\{(IH)\theta _i\}`$ converges in $`W^{m+2,p}`$.
Let us denote by $`\psi =\underset{i\mathrm{}}{lim}(1H)\theta _i`$, we then have $`\omega =\mathrm{\Delta }\psi \mathrm{\Delta }\left(W^{m+2,p}(M)\right)`$.
Step 5. Let us denote by $`^{m,p}=\mathrm{ker}HW^{m,p}(M,\mathrm{\Lambda }^k)=\mathrm{Im}(IH)W^{m,p}(M,\mathrm{\Lambda }^k)`$. Then $`\mathrm{\Delta }:^{m+2,p}^{m,p}`$ is continuous, injective and has closed image by previous step. Furthermore, $`\mathrm{Im}\mathrm{\Delta }^{m,p}`$ is dense because any smooth form in $`^{m,p}`$ is the image under $`\mathrm{\Delta }`$ of a smooth form in $`^{m+2,p}`$. To sum up, we have proved that
$$\mathrm{\Delta }:^{m+2,p}^{m,p}$$
is a continuous linear bijection.
Step 6. By the Banach open mapping theorem, we finally see that
$$G=\mathrm{\Delta }^1(1H):W^{m,p}(M,\mathrm{\Lambda }^k)^{m+2,p}W^{m+2,p}(M,\mathrm{\Lambda }^k)$$
is a bounded operator. |
warning/0506/nlin0506038.html | ar5iv | text | # Anomalous Scaling of Structure Functions and Dynamic Constraints on Turbulence Simulations
## 1 Background
The theory of turbulence and the development of calculation methods for high-Reynolds-number flows became an active research topic around the beginning of the twentieth century. This effort yielded many important results of general interest in statistical physics. For instance, Kolmogorov’s work - on turbulence theory formulated the scaling ideas for the first time, and Kraichnan proposed the mode coupling approach. However, the “turbulence problem”, lacking a small parameter characterizing the strong nonlinear interactions, has turned out to be remarkably difficult—and it remains so today.
The revolutionary realization of Osborne Reynolds that turbulence theory is a subject of statistical hydrodynamics rather than classical hydrodynamics, led almost hundred years ago to various elegant and useful phenomenological models based on ideas of kinetic theory (Prandtl , Richardson , Kolmogorov ), which strongly impacted the engineering profession. These heuristic semi-empirical models, based on low-order closures of various perturbation expansions, had a somewhat limited range of success and needed adjustable parameters, often varying from flow to flow. Nevertheless, the role of these models was—and still is—so immense that one can hardly imagine processes in mechanical and chemical engineering, aerodynamics and meteorology which do not have their input.
With the advent of powerful computers, the possibility of accurate numerical simulations, directly based on the Navier-Stokes equations, became a reality. Since the introduction of spectral methods in the end of sixties -, direct numerical simulations (DNS) have become a new tool to attack the “turbulence problem”. A strategic goal of the DNS has been to complement expensive and complicated physical experiments, and their dream is to dispense with them altogether.
The computational power required for DNS is estimated on the basis of Kolmogorov’s phenomenology that describes turbulent fluctuations filling the interval of wavenumbers $`1/Lk1/\eta _K`$, where $`L`$ and $`\eta _K=LRe^{\frac{3}{4}}`$ are the integral and dissipation scales, respectively, and $`Re=u_{rms}L/\nu `$ is the Reynolds number based on $`L`$ and the root-mean-square velocity $`u_{rms}`$. If we assume that the velocity fluctuations on scales $`r<<\eta _K`$ are highly damped and cannot contribute to the inertial range dynamics, the effective number of degrees of freedom is then $`(L/\eta _K)^3=Re^{9/4}`$. This is the minimum number of grid points required in DNS for a cubic box of linear dimension $`L`$. The required number of time steps in the computation is usually proportional to the spatial grid points, so the total computational work increases as $`Re^3`$. This means that a mere doubling of the Reynolds number requires almost an order of magnitude increase of computational work.
The accuracy of numerical methods is traditionally estimated as follows. The dissipation contribution to the equation for turbulent kinetic energy is given by
$$=\nu \overline{𝐮\frac{^2𝐮}{x_i^2}}=\nu lim_{r\eta }\frac{^2}{r^2}\overline{u_i(x)u_i(x+r)}=\nu lim_{r\eta }\frac{1}{2}\frac{^2}{r^2}S_{2,0}(r)\nu ^{\frac{2}{3}}\eta ^{\xi _22},$$
where the order of magnitude estimate in the last step comes from Kolmogorov’s phenomenology. For this case, $`\xi _2=2/3`$ and we have $`\eta _K=(\frac{\nu ^3}{})^{\frac{1}{4}}`$. We then have the familiar estimate $`\eta _KLRe^{\frac{3}{4}}`$, mentioned earlier. Thus, to accurately describe the flow, one has to simply account for fluctuations on the scales $`r\eta _K`$ by choosing the computational mesh size to be
$$\mathrm{\Delta }=a\eta _KaLRe^{\frac{3}{4}},$$
(1)
where $`a=const=O(1)`$. On this mesh, the velocity derivative is defined as
$$\frac{u(x+\mathrm{\Delta })u(x)}{\mathrm{\Delta }}=\frac{u(x)}{x}+\underset{n=2}{}\frac{1}{n!}\frac{^nu(x)}{x^n}\mathrm{\Delta }^{n1}.$$
(2)
Now, in Kolmogorov’s turbulence, $`(_xu)_{rms}=\sqrt{\overline{(_xu)^2}}(\frac{Re}{u_{rms}L})^{\frac{1}{2}}=O(Re^{\frac{1}{2}})`$, and, since $`\frac{^nu(x)}{x^n}_xu(x)/\eta _K^{n1}`$, using the mesh size $`\mathrm{\Delta }`$ from the relation (1), we arrive at the estimate
$$\frac{1}{n!}(\frac{^nu(x)}{x^n})_{rms}\mathrm{\Delta }^{n1}\frac{1}{n!}(_xu)_{rms}(\frac{\mathrm{\Delta }}{\eta _K})^{n1}\frac{a^{n1}}{n!}Re^{\frac{1}{2}}.$$
(3)
The relation (3) is essentially the basis for all numerical finite difference schemes used for the DNS of turbulence . Indeed, we see that if $`a<1`$, the first-order finite difference accurately represents the velocity derivatives.
In spectral simulations of isotropic and homogeneous turbulence, one prescribes a suitable number of the Fourier modes to represent the velocity field. Usually, this number is chosen on the basis of the magnitude of the expected Kolmogorov scale $`\eta _K`$ or the largest wavenumber $`k_{max}=2\pi /\eta _K`$. In the state-of-the-art simulations ,, the cut-off is usually chosen such that $`k_{max}=\sqrt{2}N/3`$ on a grid of size $`N^3`$.
In summary, the principal elements of Kolmogorov’s phenomenology which have enabled these traditional estimates are the following: (a) the scaling exponents of the structure functions $`S_{n,0}r^{\xi _n}`$ are given by the Kolmogorov values $`\xi _n=n/3`$; (b) the mean dissipation rate $`=\nu \overline{(_iu_j)^2}`$ is constant and $`O(1)`$, as are the moments of the dissipation rate $`\overline{^n}`$ for all $`n`$; if the latter were not the case, one can define different Kolmogorov scales on the basis of different moments of $``$; and (c) the “skewness” factors $`\overline{(_xu)^n}/\overline{(_xu)^2}^{\frac{n}{2}}=O(1)`$, independent of the Reynolds number; for, if this were not so, one can again define different Kolmogorov scales through odd moments of different order.
The main point of the present paper is that there is a need to reexamine the traditional estimates in the light of modern developments in turbulent theory and experiment. We concentrate on isotropic and homogeneous turbulence but expect that the considerations hold for more general flows as well.
## 2 Results for Intermittent Turbulence
We are interested in the Navier-Stokes dynamics of incompressible fluids. In 1941, Kolomogorov derived the few exact relation of turbulence theory, presented here for an arbitrary space dimensionality $`d`$, as
$$\frac{1}{r^{d+1}}\frac{}{r}r^{d+1}S_{3,0}=(1)^d\frac{12}{d},$$
giving $`S_{3,0}=\frac{12}{d(d+2)}r`$ and $`S_{3,0}/S_{1,2}=3.`$ A dimensional generalization of this result, without however the analytical support, yields the Kolomogorov’s (normal) scaling $`\xi _n=n/3`$. Recently ,, some additional exact consequences of the Navier-Stokes equations have been derived. In combination with recent experimental results, we consider their consequences for intermittent turbulence.
a. Dissipation scale as a random field We consider the moments of velocity difference (also called structure functions). Choosing the displacement vector $`𝐫`$ parallel to the “$`x`$-axis”, we can define the structure functions $`S_{n,m}(r)=\overline{(u(𝐱+r𝐢)u(𝐱))^n(v(𝐱+r𝐢)v(𝐱))^n}\overline{(\delta _ru)^m(\delta _rv)^n}`$, where $`u`$ and $`v`$ are the components of velocity vector parallel and normal the $`x`$-axis, respectively. In the inertial range the velocity structure functions are $`Re`$-independent; that is, if the displacement $`r`$ belongs to the interval $`\eta rL`$, then $`S_{n,m}(r)`$ do not involve any information about the dissipation scale.
Modern experiments have revealed that Kolmogorov’s result $`\xi _n=n/3`$ is almost certainly incorrect and that $`\xi _n`$ is a concave function of $`n`$—or the ratio $`\xi _n/n`$ is a decreasing function of the moment number $`n`$. (See for example Refs. for reviews and Ref. for the most recent data.) Further, the form of structure functions is given by $`S_{2n}(r)=\overline{(u(x+r)u(x))^{2n}}(2n1)!!(ϵL)^{\frac{2n}{3}}(\frac{r}{L})^{\xi _{2n}}`$. The factor $`(2n1)!!`$, ensuring Gaussian statistics at the integral scale $`L`$, is a subject of a forthcoming paper, but it suffices here to say here that it has been recently verified in experiments and numerical simulations . On the other hand, in the limit $`r0`$, the analytic structure function is approximately equal to $`S_{2n}(r)\overline{(_xu(0))^{2n}}r^{2n}`$. Combining the two, we can define a natural dissipation scale of the $`2n^{th}`$-order structure function - as
$$\eta _{2n}=(\overline{(_xu)^{2n}})^{\frac{1}{\xi _{2n}2n}}((2n1)!!ϵ^{\frac{2n}{3}}L^{\frac{2n}{3}\xi _n})^{\frac{1}{2n\xi _{2n}}}.$$
(4)
According to (4), the dissipation scales, which are expressed in terms of the moments of velocity derivatives, define a random field $`\eta `$. By a random field we mean here that the value of the length scale $`\eta `$ depends on the order of the moment considered. It will be shown below that (4) is an approximation to a more accurate representation. Similar ideas were proposed earlier in Refs. - within the framework of multifractal theories. Writing $`i_{2n}=[(2n1)!!]^{\frac{1}{2n\xi _{2n}}}`$, and using the Stirling formula $`(n1)`$, one obtains $`i_{2n}(\frac{n}{2e})^{\frac{3}{4}}`$ for $`\xi _n=n/3`$. This means that the effect of the factor $`(2n1)!!`$ can be safely neglected. For anomalous exponents $`\xi _n<n/3`$, this factor is even closer to unity and does not modify the conclusions obtained below.
b. Dissipation anomaly If the velocity field is differentiable, we obtain $`S_3(r)r^3`$ and $`_rS_3(r)0`$ in contradiction with the Kolmogorov relation. This implies that the velocity field is singular in the limit of $`\nu 0`$ and $`r0`$ (in that order), leading to the so-called dissipation anomaly. Here we first reproduce some details of Polyakov’s derivation of the dissipation anomaly for turbulence governed by Burgers equation and then outline similar procedure for the Navier-Stokes equations. Consider the one-dimensional Burgers equation
$$\frac{u}{t}+u\frac{u}{x}=\nu \frac{^2u}{x^2},$$
(5)
for which the energy balance reads as
$$\frac{1}{2}\frac{u^2}{t}+\frac{1}{3}\frac{}{x}u^3=\nu u(x)\frac{^2u}{x^2}.$$
Introducing $`x_\pm =x\pm \frac{y}{2}`$, so that, $`\frac{1}{2}\frac{}{x_\pm }=\pm \frac{}{y}`$, we can represent the energy balance equation as
$$lim_{y0}[\frac{u(x_+)u(x_{})}{t}+\frac{1}{2}\frac{}{x_+}u(x_+)^2u(x_{})+\frac{1}{2}\frac{}{x_{}}u(x_{})^2u(x_+)=\nu (\frac{^2}{x_+^2}+\frac{^2}{x_{}^2})u(x_+)u(x_{}))].$$
(6)
We also have the identities:
$$\frac{}{y}(u(x_+)u(x_{}))^3=\frac{1}{2}[\frac{u(x_+)^3}{x_+}+\frac{u(x_{})^3}{x_{}}]\frac{3}{2}[\frac{u(x_+)^2u(x_{})}{x_+}+\frac{u(x_{})^2u(x_+)}{x_{}}],$$
(7)
and
$$\nu [u(x_+)\frac{^2u(x_{})}{x_{}^2}+u(x_{})\frac{^2u(x_+)}{x_+^2}]=\nu [(u(x_+)u(x_{}))\frac{^2}{y^2}(u(x_+)u(x_{}))]+D,$$
(8)
where $`D`$, the dissipation contribution to the energy balance, is given by
$$D=\nu [u(x_+)\frac{^2}{x_+^2}u(x_+)+u(x_{})\frac{^2}{x_{}^2}u(x_{})].$$
(9)
Substituting these identities into the equation (6) and taking account of the fact that $`lim_{y0}\frac{u(x_\pm )^3}{x_\pm }=\frac{u(x)^3}{x}`$, so that in the limit $`y0`$ all non-singular terms disappear by virtue of the energy equation (5), we are left with the balance between the singular (anomalous) contributions
$$\underset{y0}{lim}\frac{1}{6}\frac{}{y}(u(x_+)u(x_{}))^3=\nu [(u(x_+)u(x_{}))\frac{^2}{y^2}(u(x_+)u(x_{}))].$$
(10)
This is Polyakov’s expression for the dissipation anomaly derived for the Burgers equation . Averaging (10) gives the exact relation $`\overline{(\delta _yu)^3}=12y`$ where the dissipation rate $`=\nu \overline{(\frac{u}{x})^2}`$.
We are interested in the Navier-Stokes dynamics of incompressible fluids, for which the energy balance equation (with the density $`\rho =1`$) is written as
$$\frac{1}{2}\frac{u^2}{t}+\frac{1}{2}𝐮u^2=p𝐮+\nu 𝐮\frac{^2𝐮}{x_i^2},$$
and that for the scalar product $`𝐮(𝐱+\frac{𝐲}{\mathrm{𝟐}})𝐮(𝐱\frac{𝐲}{\mathrm{𝟐}})𝐮(+)𝐮()`$ can be written as
$`{\displaystyle \frac{𝐮(+)𝐮()}{t}}+𝐮(+){\displaystyle \frac{}{𝐱_+}}𝐮(+)𝐮()+𝐮(){\displaystyle \frac{}{𝐱_{}}}𝐮()𝐮(+)=`$
$`{\displaystyle \frac{p(+)}{x_{+,i}}}u_i(){\displaystyle \frac{p()}{x_{,i}}}u_i(+)+\nu [𝐮(){\displaystyle \frac{^2}{x_{+,j}^2}}𝐮(+)+𝐮(+){\displaystyle \frac{^2}{x_{,j}^2}}𝐮()].`$ (11)
It is clear that in the limit $`y0`$, for which $`𝐱_\pm 𝐱`$, this equation gives the energy balance. Following Polyakov’s procedure outlined above, let us consider the two identities:
$`{\displaystyle \frac{}{y_i}}(u_i(+)u_i())(u_j(+)u_j())^2=`$
$`{\displaystyle \frac{1}{2}}{\displaystyle \frac{}{x_{+,i}}}u_i(+)u_j^2(+)+{\displaystyle \frac{1}{2}}{\displaystyle \frac{}{x_{+,i}}}u_i(+)u_j^2(){\displaystyle \frac{}{x_{+,i}}}u_i(+)u_j(+)u_j()+`$
$`{\displaystyle \frac{1}{2}}{\displaystyle \frac{}{x_{,i}}}u_i()u_j^2()+{\displaystyle \frac{1}{2}}{\displaystyle \frac{}{x_{,i}}}u_i()u_j^2(+){\displaystyle \frac{}{x_{,i}}}u_i(+)u_j()u_j(+)`$ (12)
and
$`u_i(+){\displaystyle \frac{^2}{x_{,j}^2}}u_i()+u_i(){\displaystyle \frac{^2}{x_{+,j}^2}}u_i(+)=`$
$`4(u_i(+)u_i()){\displaystyle \frac{^2}{y_j^2}}(u_i(+)u_i())+u_i(+){\displaystyle \frac{^2}{x_{+,j}^2}}u_i(+)+`$
$`u_i(){\displaystyle \frac{^2}{x_{,j}^2}}u_i().`$ (13)
Similar identities for the pressure terms can be written easily. Substituting them into (11) and, as in the case of Burgers equation considered above, accounting for the energy balance, one has
$`lim_{y0}[{\displaystyle \frac{}{y_i}}(u_i(+)u_i())(u_j(+)u_j())^2+{\displaystyle \frac{1}{2}}({\displaystyle \frac{}{x_{+,i}}}u_i(+)u_j()^2+{\displaystyle \frac{}{x_{,i}}}u_i()u_j(+)^2)=`$
$`4\nu (u_i(+)u_i()){\displaystyle \frac{^2}{y_j^2}}(u_i(+)u_i())+({\displaystyle \frac{p(+)}{𝐱_+}}{\displaystyle \frac{p()}{𝐱_{}}})(𝐮(+)𝐮())].`$ (14)
This equation can be written in a compact form as
$$lim_{y0}[\frac{}{y_i}\delta u_i|\delta _𝐲𝐮|^2+\frac{1}{2}(\frac{}{x_{+,i}}u_i(+)u_j()^2+\frac{}{x_{,i}}u_i()u_j(+)^2)=2\delta _𝐲𝐮\delta _𝐲𝐚],$$
where $`𝐚=p+\nu ^2𝐮`$ is the Lagrangian acceleration. The equation (14) is exact. Choosing the displacement vector along one of the coordinate axes and averaging (14), one obtains
$$\frac{}{y}\overline{\delta u|\delta 𝐮|^2}=8\overline{\delta u_i\frac{^2}{y^2}\delta u_i}=2\overline{(\delta _yu_i)_x^2(\delta _yu_i})=\frac{4}{3},$$
where $`\delta _yu=\delta _𝐲𝐮𝐲/y`$. The pressure terms in and the second contribution to the left side of (14) disappeared by the averaging procedure. In general, we can choose a sphere of radius $`y<<R0`$ around a point $`𝐱`$ and average (14) over this sphere. This causes the all scalar-velocity contributions to (14) disappear and the resulting equation can be perceived as a local form of the $`4/3`$ Kolmogorov law. This fact has been realized before. Introducing the angular averaging, Robert and Duchon and Eyink locally expressed the relation (14) in terms of longitudinal and transverse velocity differences. We are interested in the order of magnitude estimates (see below), and restrict ourselves to (14).
c. Relations between the moments In the isotropic and homogeneous turbulence, the Navier-Stokes equations lead to the following exact relations for structure functions. They were derived in and and experimentally investigated in some detail in Ref. ; see also Ref. . The relations for different values of $`n`$ are
$$\frac{S_{2n,0}}{r}+\frac{d1}{r}S_{2n,0}=\frac{(2n1)(d1)}{r}S_{2n2,2}+(2n1)\overline{\delta _ra_x(x)(\delta _ru)^{2n2}}.$$
(15)
Similar equations for all structure functions $`S_{n,m}`$ can easily be obtained from the equation for generating functions derived in .
d. The closure problem Equation (15), which includes both velocity and Lagrangian acceleration increments, is not closed and cannot be solved unless the relation between acceleration and velocity differences is established. It has been proposed in Ref. that the local expression (14) written for the displacement magnitudes corresponding to the bottom of inertial range, i.e. in the limit $`y\eta 0`$ can be used as a closure. At the present time, this can be done only approximately. Since at the values of displacement $`y\eta 0`$, the difference $`\delta _yu\frac{u(0)}{x}y`$, we can modify the $`lim`$ operation in (14) as
$$\underset{y0}{lim}\underset{y\eta 0}{lim},$$
(16)
leading to the order-of-magnitude estimate
$$lim_{y\eta }A\frac{(\delta _yu)^3}{y}+B\frac{}{y}\delta _yu(\delta _yv)^2\nu \delta _yu\frac{^2}{y^2}\delta _yu\frac{\delta _yp(x)}{y}\delta _yu\delta _\eta u\delta _\eta a_x,$$
(17)
where $`A`$ and $`B`$ are undetermined constants. On extrapolating to the dissipation scale $`\eta `$ where all terms in the right side of (18) are of the same order, we derive the estimate as
$$\nu \eta \delta _\eta u\eta (u(x+\eta )u(x)).$$
(18)
The relation (18) tells us that each velocity fluctuation $`\delta _\eta u`$ is dissipated on its ‘own’ dissipation scale $`\eta `$ and the local value of the Reynolds number $`Re_l=O(1)`$. This allows a simple physical interpretation that the dissipation processes at all levels $`n`$ happen on “quasi-laminar structures” where the inertial and viscous terms are of the same order. In general, the higher the moment order, the more the intense events contribute, and the smaller the value of the corresponding dissipation scale.
e. Dissipation scales and moments of derivatives The theory gives for the moments of Lagrangian acceleration $`𝐚=p+\nu ^2𝐮`$ the result that
$$a_x\frac{\delta _\eta u}{\tau _\eta }\frac{(\delta _\eta u)^2}{\eta }\frac{(\delta _\eta u)^3}{\nu }=(\delta _\eta u)^3\frac{Re}{u_{rms}L},$$
(19)
where the turn-over time $`\tau _\eta \eta /\delta _\eta u`$.
Below we will mainly discuss the equations for even-order structure functions, for which, if the displacement $`r`$ is in the inertial range, the dissipation contribution to the increment of Lagrangian acceleration is negligibly small ,. For this case, we have
$$\frac{S_{2n,0}}{r}+\frac{d1}{r}S_{2n,0}=\frac{(2n1)(d1)}{r}S_{2n2,2}(2n1)\overline{\delta _rp_x(\delta _ru)^{2n2}},$$
(20)
where $`p_x=_xp(x)`$ and $`d`$ denotes, as before, the space dimensionality.
The relation (15) is valid for all magnitudes of displacement $`rL`$, including $`r\eta `$. Below, to simplify the notation, we will omit the subscript $`x`$ in the $`x`$-component of acceleration $`a_x`$. In this limit, treating (19) as $`a=lim_{r\eta }(\delta _ru)^3/\nu `$ and substituting it in (15) gives $`\frac{S_{2n}(r)}{r}\frac{S_{2n+1}(r)}{\nu }`$. On a scale $`r=\eta _{2n}`$, writing $`S_{n,0}A_n\eta _n^{\xi _n}`$, equation (15) gives
$$\eta _nLRe^{\frac{1}{\xi _n\xi _{n+1}1}}.$$
(21)
For Kolmogorov turbulence with $`\xi _n=n/3`$ the formula (21) reads, as expected, as $`\eta _n\eta _K=LRe^{\frac{3}{4}}`$ which is $`n`$-independent. In intermittent turbulence, where the exponents can be well-described , by the relation $`\xi _n0.383n/(1+0.05n)`$, the relation (21) defines the Reynolds-number-dependent dissipation scales. As $`n\mathrm{}`$, $`\eta _nLRe^1`$. Thus, to resolve all fluctuations including the strongest, the computational work need to increase as $`Re^4`$, as already noted in Ref. . In general, in the limit $`n\mathrm{}`$, the relation (21) can be written as
$$\eta _nLRe^{\frac{1}{\frac{d\xi _n}{dn}+1}},$$
so one may get a somewhat different estimate for the computational work than $`Re^4`$, but the principal conclusion is inescapable that intermittency makes DNS more expensive than previously thought.
Using the relations (18), (20) and (21), obtained by balancing various terms in the exact dynamic equations (14), (15), we can develop the multi-scaling algebra. For example,
$$\overline{a^{2n}}(\frac{Re}{u_{rms}L})^{2n}S_{6n}(\eta _{6n})(\frac{Re}{u_{rms}L})^{2n}\eta _{6n}^{\xi _{6n}}(\frac{u_{rms}^2}{L})^{2n}Re^{a_{2n}},$$
(22)
with $`a_{2n}=2n+\frac{\xi _{6n}}{\xi _{6n}\xi _{6n+1}1}`$. With $`\xi _6=2`$ and $`\xi _7=7/3`$, we recover Yaglom’s result $`\overline{a^2}\frac{u_{rms}^{\frac{9}{2}}}{\sqrt{\nu }}`$. The intermittency corrections are readily found from (22). Recent experiments by Reynolds et al. have lent strong support to this result. The formula (22) shows that the second moment of Lagrangian acceleration is expressed in terms of the sixth-order structure function evaluated on its dissipation scale $`\eta _6`$. To extract information about the fourth moment $`\overline{a^4}`$, we should have accurate data on $`S_{12}(\eta _{12})`$ which is very difficult to obtain in high-Reynolds-number flows.
The moments of velocity derivatives are evaluated easily. In accordance with (18), we have
$$\overline{(_xu)^{2n}}\overline{(\frac{\delta _\eta u}{\eta })^{2n}}\overline{(\frac{(\delta _\eta u)^2}{\nu })^{2n}}Re^{d_{2n}},$$
(23)
where $`d_{2n}=2n+\frac{\xi _{4n}}{\xi _{4n}\xi _{4n+1}1}`$.
It is important to stress that the first equality in (23) involves the averaging over two random fields $`u`$ and $`\eta `$. To perform this averaging, we have to either know the joint probability $`p(u,\eta ,r)`$ or use the functional relation between the fields given by (18). This leads to the second equation in (23) and the final result. Since $`\overline{(_xu)^2}Re`$, the relation (3) leads to a new relation between exponents
$$2\xi _4=\xi _5+1$$
which agrees extremely well with experimental data. The relation (23) differs from proposals reviewed in Ref. .
f. The role of the fluctuations of the dissipation scale Let us reexamine the relation (4). In the limit $`r0`$, the velocity field is analytic and can be expanded by Taylor series so that $`\frac{u}{x}\delta _ru/r`$. This gives $`\overline{(\frac{u}{x}r)^{2n}}S_{2n}(r)`$. When $`r\eta 0`$, we have to evaluate the mean of the ratio $`\overline{(\delta _\eta u/\eta )^{2n}}`$ which is not a trivial task, since we are dealing here with the ratio of two random fields—unless the relation (18), which expresses the dissipation scale in terms of velocity field, is used. If, however, we incorrectly assume that the dissipation scale fluctuations are independent of those of the velocity field and neglect the step leading to the last equations in the right hand side of (23), it is possible to write the moments of velocity derivative as
$$\overline{(_xu)^{2n}}\overline{(\frac{\delta _\eta u}{\eta })^{2n}}S_{2n}(\eta _{2n})/\eta _{2n}^{2n}Re^{p_{2n}},$$
(24)
where $`p_{2n}=\frac{\xi _{2n}2n}{\xi _{2n}\xi _{2n+1}1}`$. Equating expressions (23) and (24), we have
$$\frac{\xi _{2n}2n}{\xi _{2n}\xi _{2n+1}1}=2n+\frac{\xi _{4n}}{\xi _{4n}\xi _{4n+1}1},$$
(25)
subject to the constraints $`\xi _0=0`$ and $`\xi _3=1`$. The only solution to (25) is $`\xi _n=n/3`$. Since equation (25) is based on the first equality (23), which in general is incorrect, we can conclude that the source of anomalous scaling in hydrodynamic turbulence is the fluctuation of the dissipation scale field $`\eta `$, which itself is strongly correlated the velocity field fluctuations via expression (18). This does not preclude a different situation from arising in other forms of turbulence, e.g., scalar turbulence generated by white-noise forcing .
It follows that $`\overline{(\frac{u}{x})^2}=lim_{r\eta _2}\overline{\frac{u(x)}{x}\frac{u(x^{})}{x^{}}}=lim_{r\eta _2}\frac{^2}{r^2}\overline{u(x)u(x^{})}(2\xi _2)\eta _2^{\xi _22}.`$ The higher-order derivatives are evaluated in a similar way to yield
$$(\frac{^nu}{x^n})_{rms}=lim_{r\eta _2}\sqrt{\frac{^{2n}}{r^{2n}}S_2(r)}\eta _2^{\frac{\xi _22n}{2}}Re^{\frac{\xi _22n}{2(\xi _22)}}=Re^{\frac{1}{2}}Re^{\frac{n1}{\xi _22}}.$$
(26)
## 3 Implications for Numerical Methods
According to experimental data (see Refs. for recent results), the exponent $`\xi _20.700.71>2/3`$ and as $`n\mathrm{}`$, the terms in the expansion (2) for simulating the “typical” velocity derivatives can be estimated via
$$(\frac{^nu}{x^n})_{rms}\mathrm{\Delta }^{n1}Re^{\frac{1}{2}}Re^{\gamma (n1)},$$
(27)
with $`\gamma =(\frac{3}{4}\frac{1}{\xi _22}))>0`$. For $`\xi _20.71`$, we find $`\gamma 0.025`$. The accuracy of the numerical method in calculating the most intense velocity fluctuations can be estimated if, in the limit $`n\mathrm{}`$, the expression
$$\overline{(\frac{u}{x})^{2n}}^{\frac{1}{2n}}(\frac{\mathrm{\Delta }}{\eta _{2n}})^{n1}Re^{\frac{1}{2}}Re^{\frac{n+1}{4}}$$
(28)
is used instead of $`(_xu)_{rms}`$. In the above equation, the mesh size $`\mathrm{\Delta }`$ is defined by (1) and the expressions (23) for the moments of velocity derivative have been used. We see that when the Reynolds number is large, the high-order derivatives in the expression (2) dominate. This means that the DNS based on the mesh equal to the Kolmogorov scale becomes quite inaccurate. It is easy to check that accurate simulations of the largest fluctuations requires the resolution of the smallest scales which are $`O(1/Re)`$. This means that the computational resolution scales as $`Re^3`$ and the computational work grows as $`Re^4`$.
In Refs. , it has argued that the intermittent nature of turbulence makes the size of the attractor smaller than the conventionally estimated, so the computational power needed becomes correspondingly smaller than the conventional estimate—not larger as just claimed. The rationale is roughly that the “interesting” parts of the flow occupy small volumes of space so any reasonable computational effort that focuses on those volumes is likely to be less expensive. This is also the spirit of adaptive meshing . Even if the interesting parts of a turbulent flow are not space-filling, as discussed at length in Ref. , we do not yet know how to track them efficiently in hydrodynamics turbulence. We also do not know if the part of the flow that contains the less interesting parts can be computed with greater economy. Nevertheless, it must be said that the present estimates apply to uniform meshing, which has been the most successful of the computing schemes until now. It should also be mentioned that there is a specific suggestion on the most singular structure in turbulence, which yields $`Re^{3.6}`$, which is slightly different from $`Re^4`$ estimated in this paper.
## 4 Dynamic Constraints on Sub-Grid Models for LES
If the Reynolds number is large, the computational work involved in the numerical simulation of a flow is huge. It is interesting that at about the same time that DNS came into being, the idea of the Large Eddy Simulations (LES) was proposed by Deardorff . The idea is very simple. Consider the Navier-Stokes equations
$$_t𝐮+u_i_i𝐮=p+\nu ^2𝐮;_iu_i=0.$$
(29)
We choose the mesh size $`\mathrm{\Delta }`$ and define the so-called “sub-grid” velocity fluctuations $`u^>(k)0`$ for $`k\pi /\mathrm{\Delta }`$. The Fourier-transform of velocity field is defined as
$$u(𝐤)=u^<(𝐤)+u^>(𝐤),$$
(30)
so that
$$u^>(𝐱)=_{|k|>\frac{2\pi }{\mathrm{\Delta }}}e^{i𝐤𝐱}u^>(𝐤)d^3k;u^<(𝐱)=_{|k|\frac{2\pi }{\mathrm{\Delta }}}e^{i𝐤𝐱}u^<(𝐤)d^3k.$$
(31)
The goal is to obtain the correct equation for the resolved scales $`u^<(k)0`$ in the interval $`0k\pi /\mathrm{\Delta }`$. We decompose the field and write the equation for only the resolved scales as
$$_t𝐮^<+u_i^<_i𝐮^<=𝒮𝒢p^<+\nu ^2𝐮^<,$$
(32)
where, for this particular formulation, the subgrid contribution is $`𝒮𝒢=u_i^<_i𝐮^>u_i^>_i𝐮^<u_i^>_i𝐮^>`$. The LES equations are considered a success if the large-scale velocity fields (i.e., for $`k1/\mathrm{\Delta }`$) given by the Navier-Stokes equations (30) and by a model (33) are identical or close enough for all Reynolds numbers.
There is, however, one problem. To derive the equation of motion containing only the resolved fields, one has to express all contributions to $`𝒮𝒢`$, involving the sub-grid velocity fluctuations $`𝐮^>`$, in terms of $`𝐮^<`$, which is basically equivalent to solution of the proverbial “turbulence problem”. The model equation (33) is written in a generic form, but a similar difficulty arises if, instead of the Fourier-space decomposition introduced above, the filtering or any other kind is used.
The accurate LES model must satisfy the following dynamic constraints. The method developed in the Ref. can be literally applied to the Navier-Stokes equations with an arbitrary right hand side and, defining the coarse-grained structure functions $`S_{n,0}^<(r)=\overline{(\delta _ru^<)^n}`$, we obtain, from (21), the result
$$\frac{S_{2n,0}^<}{r}+\frac{d1}{r}S_{2n,0}^<=\frac{(2n1)(d1)}{r}S_{2n2,2}^<+(2n1)\overline{(\delta _r(𝒮𝒢_x)\delta _rp_x^<)(\delta _ru^<)^{2n2}}.$$
(33)
The large-scale velocity fields obtained from DNS and LES can be identical $`S_{n,0}(r)=S_{n,0}^<(r)`$ if and only if
$$\overline{(\delta _r(𝒮𝒢_x)\delta _rp_x^<)(\delta _ru^<)^{2n2}}=\overline{\delta _rp_x(\delta _ru)^{2n2}}.$$
(34)
Similar constraints, coming from the equations for various structure functions $`S_{n,m}`$ can be readily obtained. It is impossible to demand equality of two random fields $`𝐮`$ and $`𝐮^<`$ obtained from two different equations. The only criterion we can impose is that of statistical equality or, equivalently, constraint on all moments, namely $`S_n^<(r)=S_n(r)`$. The relations (34), reflecting this necessary condition of the LES validity, must be satisfied.
We wish to stress that these constraints are not dissimilar to $`S_{n,m}^{LES}S_{n,m}^<`$, often implied in the literature. Here $`S_{n,m}^{LES}(r)`$ are the structure functions evaluated from the velocity field obtained from LES. The velocity increment can be written as $`\delta _ru=u(k)e^{ikx}(e^{ikr}1)`$, so that
$$S_2E(k)(1\mathrm{cos}kr)𝑑k.$$
It is easy to see that if $`r<<L`$, where $`L`$ is the integral scale, and the energy spectrum decreases with $`k`$ fast enough, the main contribution to the integral comes from the range where $`kr1`$. Thus the structure functions $`S_{n,0}(r)`$ probe structures on the scales of the order $`r`$ and cannot differ strongly from the one obtained from the filtered field.
Various model considerations, leading to expressions for $`𝒮𝒢`$, have been suggested in the last forty years. Consider the example that follows from Kolmogorov’s theory. If the role of the small scale fluctuations in the large-scale dynamics can be expressed in terms of effective viscosity $`\nu _{SG}`$, then $`\nu _{SG}(\overline{ϵ}\mathrm{\Delta }^4)^{\frac{1}{3}}`$. Then, dropping the averaging sign (quite an assumption!) and substituting a simple estimate coming from the energy balance, namely, $`ϵ=\nu _{SG}S_{ij}^<S_{ij}^<\nu _{SG}S_{ij}^2`$, we derive the Smagorinsky formula given by $`\nu _{SG}=\alpha \sqrt{S_{ij}^<S_{ij}^<}\mathrm{\Delta }^2`$, where $`\alpha =O(1)`$. It is important that the resolved rate of strain is evaluated in terms of velocity differences on the computational mesh
$$S_{ij}^<(x)=\frac{1}{2}(\frac{u_i^<(x+\mathrm{\Delta }_j)u_i^<(x)}{\mathrm{\Delta }_j}+\frac{u_j^<(x+\mathrm{\Delta }_i)u_j^<(x)}{\mathrm{\Delta }_i}),$$
(35)
where $`i,j=1,2,3`$. In this approximation, the Reynolds stress $`\tau _{ij}=\overline{u_iu_j}\nu S_{ij}\nu _{SG}S_{ij}^<`$. Equation (36) with the model for $`𝒮𝒢`$ defines a closed set of equations which can be used for LES. The analytically evaluated coefficient from Yakhot and Orszag gives $`\alpha 0.2`$, while the so-called dynamic method gives something different. In all approaches, since the large-scale fields $`\delta _r𝐮^<`$ and $`\delta _r𝐮`$ are statistically independent upon Reynolds number, the parameter $`\alpha =O(Re^0)`$. Thus, this simple model is
$$𝒮𝒢a\mathrm{\Delta }^2|S_{ij}^<|u^<=O(1).$$
(36)
Examining the relations (34) and (36), an interesting conclusion can be reached. If $`\mathrm{\Delta }r`$, one can assume statistical independence of all velocity differences $`\delta _ru^<`$ and $`\delta _\mathrm{\Delta }u^<`$. Since $`𝒮𝒢`$ given by (34) and (35) depends on the velocity differences defined on the mesh size $`\mathrm{\Delta }`$ as
$$\overline{\delta _r𝒮𝒢(\delta _ru^<)^{2n2}}\overline{\delta _r𝒮𝒢}\overline{(\delta _ru^<)^{2n2}}=0,$$
(37)
we see that the Smagorinsky model satisfies the dynamic constraints, provided the pressure gradient differences in the filtered and unfiltered fields are close to each other. The validity of the dynamic Smagorinsky models in the range $`k<<1/\mathrm{\Delta }`$ has been verified by large eddy simulations (A. Oberai, private communication 2005). However, as $`r\mathrm{\Delta }`$, $`\delta _r𝒮𝒢`$, $`\delta _rp_x`$ and $`\delta _ru^<`$ are strongly correlated and, as a result, the model becomes invalid. This consideration is applicable to all low-order closures.
This intrinsic failure of all existing LES models at scales comparable to the computational mesh is well-known. At sufficiently low Reynolds numbers, LES give accurate results. However, with increase of $`Re`$ the quality of the simulations deteriorates starting from the vicinity of the cut-off, propagating toward the larger scales. At this point one is forced to increase the resolution. The reasons for this failure can be qualitatively understood as follows. Consider LES at a relatively low Re on a fixed mesh $`\mathrm{\Delta }/L_1=\gamma _1`$ where $`L_1`$ is an integral scale of this particular simulation. Now increase the length scale of the flow $`L_2>>L_1`$, thus increasing the Reynolds number. If, in the first case, the number of the cascade steps for the energy flux to reach the mesh scale was say $`n_1`$, that in the second simulation is equal to $`n_2>>n_1`$. Since the intermittency and deviation from the close-to-Gaussian statistics, experimentally observed at the integral scale, grows with the number of cascade steps, the contribution from the very strong velocity fluctuations at the “dissipation” scale $`\mathrm{\Delta }`$ increases. As a result, the low order models that are successful in the close to Gaussian situations break down. In another scenario, let us increase the Reynolds number by increasing the mean velocity while keeping both the energy injection scale and the mesh size $`\mathrm{\Delta }`$ constant. In this situation, the top of the “inertial” range will move into the range of scales which are larger than $`\mathrm{\Delta }`$, thus again invalidating the LES.
A recent paper by Kang et al. has demonstrated that, at the scales close to those of the mesh size, the probability density function $`p(\delta _ru)`$ computed from LES was quite close to a Gaussian while the experimental PDF showed broader tails, typical of intermittency. This means that the contributions from strong velocity fluctuations obtained from LES are underpredicted. Since the intermittent effects becomes stronger with increasing Reynolds number, we expect this difference to grow, thus invalidating the LES if the mesh size is also not modified. A very interesting example is given by the LES of the flow in a simple cavity reported by Larcheveque et al. . It was shown that to correctly reproduce the experimental data on pressure fluctuations in a frequency range $`100f2000Hz`$, the optimal cut-off of the large eddy simulations corresponded to $`\mathrm{\Delta }_f=100KHz`$. With decrease of $`\mathrm{\Delta }_f`$, the quality of the results rapidly deteriorated. The present theory explains the failure of LES schemes with fixed mesh to describe the high Reynolds number flows as originating from the failure of low-order models in an all-important range $`r\mathrm{\Delta }`$, this range being responsible for the energy cascade dissipation. At the present time, it is not clear how many constraints (35) must be satisfied to achieve accurate LES, but we believe that the number must grow with the Reynolds number.
## 5 Conclusions
For many years, intermittency and anomalous scaling in three-dimensional turbulence were considered major challenges for theorists. Recent developments of the multifractal theory and its dynamic formulation led to description of intermittency in terms of an infinite number of dissipation scales (ultraviolet cut-offs). It was shown that strong velocity fluctuations are dissipated on scales that are much smaller than that estimated from Kolomogorov’s theory. In this paper, we have attempted to make a connection between the theory of anomalous scaling and numerical methods.
One conclusion that follows from this connection is that to simulate all fluctuations, including the strongest ones, the computational demands scale as $`Re^4`$, and not as $`Re^3`$ as traditionally deduced according to the Kolmogorov theory. To achieve the full DNS of turbulence, including the strongest small-scale velocity fluctuations, one has to use resolutions high enough to produce an analytic interval of structure functions, where $`S_n\overline{_xu(0))^n}r^n`$. Analyzing the results of various numerical state-of-the-art DNS, we have discovered that this criterion is satisfied only for $`n4`$. This is not sufficient to accurately simulate the velocity derivatives.
A second comment concerns the Large Eddy Simulations. An infinite number of dynamic constraints on a correct subgrid model has been derived from the exact relations for structure functions. Due to the Galilean invariance, the subgrid scales cannot influence the advective term in the Navier-Stokes equations, provided the subgrid scale $`\mathrm{\Delta }/r0`$. However, it is clear from analyzing the equations of Section 4 that the subgrid model cannot be reduced to a low-order viscosity expression, but must include high-order nonlinear contributions that do not vanish at the scales close to the mesh size.
Thus, while accurate DNS are possible if the resolution requirements are met and powerful enough computers are available. However, due to the basic theoretical problems, derivation of an accurate and theoretically justified subgrid model, valid at very high Reynolds numbers, remains a major challenge.
It is worth pointing out that we have considered homogeneous and isotropic turbulence. The situation with wall flows is even more complex. There, turbulence is mainly produced in the vicinity of the wall where acceleration and turbulence production are highly intermittent. Recent DNS by Lee et al. have demonstrated strong intermittency and the Reynolds number dependence of the few first moments of Lagrangian acceleration near the wall, sharply peaking at the reduced normalized distance $`y_+2.5`$. At present, we do not know how to model this near-wall phenomenon that is largely responsible for turbulence production.
We wish to conclude on a “positive” note. The fact that the structure functions $`S_{2n}(2n1)!!(\frac{r}{L})^{\xi _{2n}}`$ means that the velocity distribution is close to the Gaussian near $`r=L`$, and the intermittency is weak or nonexistent. It follows that simple, semi-qualitative resummations of the expansions in powers of the dimensionless rate-of-strain are much less problematic there. Thus, the derivation of the VLES or time-dependent RANS appears to have a brighter future.
References
1. A.N. Kolmogorov, Dokl. Akad. Nauk SSSR. 30, 9 (1941)
2. A.N. Kolmogorov, Dokl. Akad. Nauk SSSR. 32, 16 (1941)
3. A.N. Kolmogorov, Izv. Akad. Nauk SSSR. Ser. Fiz. $`\mathrm{𝐈𝐕}`$ (1-2), 56 (1942)
4. R.H. Kraichnan, J. Fluid Mech. 5, 497 (1959)
5. L. Prandtl, Math. Mech. 1, 431 (1925)
6. L.F. Richardson, Proc. Roy. Soc. A10, 709 (1926)
7. S.A. Orszag and G.S. Patterson, “Numerical Simulations of Turbulence”, in Statistical Models and Turbulence, Lecture Notes in Physics 12, 127 (1972), eds. M. Rosenblatt and C.W. Van Atta, Springer, Berlin
8. S.A. Orszag, in “Numerical Simulations of the Taylor-Green vortex”, in Computing Methods in Science and Engineering. Part 2. Lecture Notes in Computer Science 11, 54 (1974)
9. The traditional estimate is due to L. Landau and E.M. Lifschitz, Fluid Mechanics, Addison Wesley, Reading, MA, 1959; it was commented upon extensively in the context of spectral methods by S.A. Orszag, “Statistical theory of turbulence”, in Fluid Dynamics, Les Houches, 2237 (1973), eds. R. Balian and J.L. Peube, Gordon and Breach, New York
10. T. Gotoh, D. Fukayama and T. Nakano, Phys. Fluids 14, 1065 (2002); Y. Kaneda, T. Ishihara, M. Yokokawa, K. Itakura and A. Uno, Phys. Fluids 15, L21 (2003)
11. P.K. Yeung, D.A. Donzis, Phys. Fluids xx, Lxx (2005)
12. V. Yakhot, Phys. Rev. E 63, 026307 (2001)
13. R.J. Hill, J. Fluid Mech. 368, 317 (2002)
14. U. Frisch, Turbulence. Legacy of A.N. Kolmogorov, Cambridge University Press, England (1995); K.R. Sreenivasan and R.A. Antonia, Annu. Rev. Fluid Mech. 29, 435 (1997)
15. S.Y. Chen, B. Dhruva, S. Kurien and M.A. Taylor, J. Fluid. Mech. 533, 183, 2005
16. Professor T. Gotoh kindly tested this relation using the experimental data published in T. Gotoh and T. Nakano, J. Stat. Phys. 113, 855 (2003); extensive tests using other sources of data have since been completed to confirm this result.
17. V. Yakhot, J. Fluid Mech. 495, 135 (2003)
18. V. Yakhot and K.R. Sreenivasan, Physica A 343, 147-155, 2004
19. G. Paladin and A. Vulpiani, Phys. Rev. A 35, 1971 (1987); G. Paladin and A. Vulpiani, Phys. Rep. 156, 147 (1990)
20. K.R. Sreenivasan and C. Meneveau, Phys. Rev. A 38, 6287 (1988); C. Meneveau and K.R. Sreenivasan, J. Fluid Mech. 224, 429 (1991)
21. M. Nelkin, Phys. Rev. A 42, 7226 (1990); U. Frisch and M. Vergassola, Europhys. Lett. 15, 439 (1901)
22. A.M. Polyakov, Phys. Rev. E 52, 6183 (1995)
23. J. Duchon and R. Robert, Nonlinearity 13, 249 (2000)
24. G.L. Eyink, Nonlinearity 16, 137 (2003)
25. S. Kurien and K.R. Sreenivasan, Phys. Rev. E 64, 056302 (2001)
26. R.A. Antonia, M. Ould-Rouis, F. Anselmet and Y. Zhu, J. Fluid Mech. 332, 395 (1997)
27. K.R. Sreenivasan, Flow, Turbulence and Combustion 72, 115 (2004)
28. A.S. Monin and A.M. Yaglom, Statistical Fluid Mechanics, vol. 2, MIT Press, Cambridge, MA (1975)
29. A.M. Reynolds, N. Mordant, A.M. Crawford and E. Bodenschatz, New J. Phys. art. no. 58 (2005)
30. M. Chetkov and G. Falkovich, Phys. Rev. Lett. 76, 2706 (1996)
31. O. Roussel and K. Schneider, Computers and Fluids 34, 817 (2005); O.V. Vasilyev and S. Paolucci, J. Comp. Phys. 125, 498 (1996)
32. Z.-S. She and E. Leveque, Phys. Rev. Lett. 72, 336 (1994)
33. J.W. Deardorff, J. Fluid Mech. 41, 453 (1970)
34. J. Smagorinsky, Mon. Weather Rev. 91, 99, (1963)
35. V. Yakhot and S.A. Orszag, J. Sci. Comp. 1, 1 (1986)
36. M. Germano, J. Fluid Mech. 238, 325 (1992)
37. H.S. Kang, S. Chester and C. Meneveau, J. Fluid Mech. 480, 129 (2003)
38. L. Larcheveque, P. Segaut, T-H. Le and P. Comte, J. Fluid Mech. 516, 265 (2004)
39. C. Lee, K. Yeo and J-Il Choi, Phys. Rev. Lett. 92, 144502-1 (2004) |
warning/0506/astro-ph0506058.html | ar5iv | text | # Statistical characteristics of the observed Ly-𝛼 forest and the shape of the initial power spectrum
## 1 Introduction
One of the most perspective methods to study the processes responsible for the formation and evolution of the structure in the Universe is the analysis of properties of absorbers observed in spectra of the farthest quasars. The great potential of such investigations was discussed already by Oort (1981, 1984) just after Sargent et al. (1980) established the intergalactic nature of the Lyman-$`\alpha `$ forest. Indeed, the absorption lines trace the small scale distribution of hydrogen along the line of sight at redshifts $`z`$ 2 when matter is not yet strongly clustered and its observed characteristics can be more easily interpreted. The available Keck and VLT high resolution observations of the Lyman-$`\alpha `$ forest provide a reasonable database and allow one to apply statistical methods for their analysis.
The composition and spatial distribution of the observed absorbers is complicated. Thus, at large redshifts the population of rich metal systems including Ly-damped and Ly-limit systems are rare and majority of observed absorbers is associated with isolated low mass HI clouds. At low redshifts a significant number of stronger Ly-$`\alpha `$ lines and metal systems is associated with galaxies (Bergeron et al. 1992; Lanzetta et al. 1995; Tytler 1995; Le Brune et al. 1996). However as was recently shown by Penton, Stock and Shull (2000; 2002) and McLin et al. (2002), even at small redshifts some absorbers are associated with galaxy filaments while others are found within galaxy voids. These results suggest that the population of weaker absorbers dominating at higher redshifts can be associated with weaker structure elements formed by the non luminous baryonic and DM components. They also suggest that the Ly–$`\alpha `$ forest can be considered as a low mass component of the generic Large Scale Structure (LSS) which is seen in simulated and observed spatial matter distribution. This means that absorbers at high $`z`$ trace the DM structure which is qualitatively similar to the rescaled one observed at small redshifts.
In this paper we investigate a sample of $`6000`$ absorbers observed in 19 high resolution spectra of QSOs and compare their properties with the improved model of absorbers proposed in Demiański & Doroshkevich (2003 a,b, hereafter Paper I & Paper II). We assume that absorbers are dominated by long–lived gravitationally bound and partially relaxed clouds composed of both DM and baryonic components. The fact that we can observe galaxies and quasars at $`z3`$ demonstrates the existence of strong density perturbations already then. Here we show that at these redshifts there are also strong negative density perturbations of a galactic scale which can be identified with rapidly expanded underdense regions. Such regions are naturally associated with the colder absorbers.
Our model explains the self similar character of evolution of the observed Doppler parameter, HI column density and absorber separation. Such evolution implies that the mean values of these characteristics slowly vary with redshift while their probability distribution functions (PDFs) remain unchanged. This model links the observed and other physical characteristics of absorbers – such as their DM column density, size, and fraction of matter associated with absorbers – and allows us also to identify several subpopulations of absorbers with different evolutionary histories.
We treat the evolution of structure as a random process of formation and merging of Zel’dovich pancakes, their transverse expansion and/or compression and successive transformation into filaments and high density clouds. Later on the hierarchical merging of pancakes, filaments and clouds forms rich galaxy walls observed at small redshifts. Impact of these factors is clearly seen in high resolution numerical simulations of evolution of the LSS (for review see, e.g. Frenk 2002).
Theoretical expectations of our model are based on the Zel’dovich approximate theory of gravitational instability (Zel’dovich 1970; Shandarin & Zel’dovich 1988). As is well known, it correctly describes only the linear and weakly nonlinear stages of the structure formation and cannot describe later stages of evolution of structure elements. In spite of this, the statistical approach proposed in (Demiański & Doroshkevich 1999, 2004a; hereafter DD99 & DD04) nicely describes the main properties of observed and simulated LSS (Demiański et al. 2000; Doroshkevich, Tucker, Allam & Way 2004) without any smoothing or filtering procedure.
Presently various observations are used to determine the power spectrum of the initial density perturbations. Its amplitude and its shape on scales $`10h^1`$Mpc are approximately established by investigations of the microwave relic radiation (Spergel et al. 2003, 2006) and the structure of the Universe at $`z`$ 1 detected in large redshift surveys (Percival et al. 2001; Tegmark, Hamilton & Xu 2003; Verde et al. 2002, 2003) and weak lensing data (see, e.g., Hoekstra, Yee & Gladders 2002).
The shape of the initial power spectrum on scales $`10h^1`$ Mpc – $`1h^1`$Mpc can be tested at high redshifts where it is not yet strongly distorted by nonlinear evolution (Croft et al. 1998, 2002; Nusser & Haehnelt 2000; Gnedin & Hamilton 2002; Viel et al. 2004 a,b; Kim et al. 2004; McDonald et al. 2004; Seljak et al. 2004; Zaroubi et al. 2005). The method used in these papers is surprisingly universal and is successfully applied to spectra observed with both high and moderate resolution. It utilizes the measured transmitted flux only and does not require preliminary determination of column density, Doppler parameters and even discrimination of hydrogen and metal line systems (McDonald et al. 2004, Seljak et al. 2004). In spite of this, it successfully restores the CDM–like power spectrum down to scales $`1h^1`$ Mpc. Recent results on reconstruction of the initial power spectrum are summarized and discussed in many papers (see, e.g., Tegmark and Zaldarriaga 2003; Wang et al. 2003; Zaldarriaga, Scoccimorro & Hui 2003; Peiris et al. 2003; Spergel et al. 2003,2006; Tegmark et al. 2004; McDonald et al. 2004, Seljak et al. 2004).
A straightforward method of reconstruction of the initial power spectrum from the observed characteristics of absorbers was proposed and tested in Paper II. This method can be used to recover the initial power spectrum down to unprecedentedly small scale. In contrast with previous investigations (Croft et al. 1998, 2002; Nusser & Haehnelt 2000; Viel et al. 2004b; McDonald et al. 2004) we analyze the separation between adjacent absorbers and their column density rather than the flux or smoothed density field. This means that our results are not restricted by the standard factors such as the Nyquist limit, the impact of nonlinear processes, the unknown matter distribution between absorbers or their peculiar velocities. This approach successfully complements investigations of the power spectrum mentioned above.
Here we improve the analysis discussed in Paper II by using a richer observed sample and a more refined model of absorbers. We use two independent methods of determination of the initial power spectrum. The first one is based on measurements of the separation between adjacent absorbers, while the second one, proposed in Paper II, uses measurements of the column density of absorbers. Both approaches allow one to determine the spectrum down to the scale of $`510h^1`$ kpc. At scales $`(0.1510)h^1`$ Mpc our results coincide with those expected for the CDM–like power spectrum and Gaussian perturbations with the precision of $``$ 15%. However, we have found some evidence that at scales $`0.15h^1`$Mpc the initial power spectrum differs from the CDM–like one suggesting a complex inflation with generation of excess power at small scales. Such excess power accelerates the process of galaxy formation at high redshifts and can shift the epoch of reheating of the Universe to higher redshifts.
At present we have only limited information about the properties of the background gas and the UV radiation (Haardt & Madau 1996; Scott et al. 2000, 2002; Schaye et al. 2000; McDonald & Miralda–Escude 2001; McDonald et al. 2000, 2001; Theuns et al. 2002 a, b; Levshakov at al. 2003; Boksenberg, Sargent & Rauch 2003; Demiański & Doroshkevich 2004b) and therefor some numerical factors in our model remain undetermined. This means that our approach should be tested on representative numerical simulations that more accurately follow the process of formation and disruption of pancakes and filaments and provide a unified picture of the process of absorbers formation and evolution (see, e.g., Weinberg et al. 1998; Zhang et al. 1998; Davé et al. 1999; Theuns et al. 1999, 2000). But so far such simulations are performed mainly in small boxes what restricts their representativity, introduces artificial cutoffs in the power spectrum and complicates the quantitative description of structure evolution (see more detailed discussion in Gnedin & Hamilton 2002; Tegmark & Zaldarriaga 2002; Zaldarriaga Scoccimoro & Hui 2002; Seljak, McDonald & Makarov 2003; Manning 2003; Paper II). As was shown by Meiksin, Bryan & Machacek (2001), the available simulations reproduce quite well the characteristics of the flux but cannot restore other observed characteristics of the forest. This means that first of all numerical simulations should be improved (see more detailed discussion in Paper II).
Comparison of results obtained in Paper I and Paper II and in this paper demonstrates that the quality and representativity of the sample of observed absorbers are very important for the reconstruction of processes of absorbers formation and evolution. Thus, richer sample of the observed absorbers makes it possible to select and investigate several representative subsamples of absorbers with different evolutionary histories. None the less, our analysis indicates a possible deficit of weaker absorbers and pairs of close absorbers what in turn could be related to insufficient sensitivity of the process of absorbers’ identification. Thus, the number and parameters of absorbers identified for the same quasar depend upon the used identification procedure. Further progress can be achieved first of all with richer samples covering the range of redshifts at least up to $`z`$ 5 what would allow one to perform a complex investigation of the early period of structure evolution.
This paper is organized as follows. In Sec. 2 the observational databases used in our analysis are presented and statistical characteristics of the observed parameters of absorbers are obtained. The theoretical models of the structure evolution are discussed in Sec. 3. The results of statistical analysis of the model dependent parameters and the derived initial power spectrum are given in Sec. 4. Discussion and conclusions can be found in Sec. 5.
## 2 Observed characteristics of absorbers
### 2.1 Properties of the homogeneously distributed matter
In this paper we consider the spatially flat $`\mathrm{\Lambda }`$CDM model of the Universe with the Hubble parameter and mean density given by:
$$H^2(z)=H_0^2\mathrm{\Omega }_m(1+z)^3[1+\mathrm{\Omega }_\mathrm{\Lambda }/\mathrm{\Omega }_m(1+z)^3],$$
$`n_b(z)=2.410^7(1+z)^3(\mathrm{\Omega }_bh^2/0.02)\mathrm{cm}^3,`$ (1)
$$\rho _m(z)=\frac{3H_0^2}{8\pi G}\mathrm{\Omega }_m(1+z)^3,H_0=100h\mathrm{km}/\mathrm{s}/\mathrm{Mpc}.$$
Here $`\mathrm{\Omega }_m=0.3\&\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$ are the dimensionless matter density and the cosmological constant (dark energy), $`\mathrm{\Omega }_b`$ is the dimensionless mean density of baryons, and $`h=`$ 0.7 is the dimensionless Hubble constant.
Properties of the compressed gas can be suitably related to the parameters of homogeneously distributed gas, which were discussed in many papers (see, e.g., Ikeuchi & Ostriker 1986; Haardt & Madau 1996; Hui & Gnedin 1997; Scott et al. 2000; McDonald et al. 2001; Theuns et al. 2002a, b; Demiański & Doroshkevich 2004b). In this paper we consider evolution of absorbers at observed redshits $`z4`$ when weak variations of the gas entropy are determined by the interaction of the gas with the UV background. Thus, the expected intensity of the UV background radiation, can be fitted as follows:
$`J(z,\nu )=J_{21}(z)\left({\displaystyle \frac{\nu _H}{\nu }}\right)^{\alpha _\gamma }10^{21}{\displaystyle \frac{erg}{scm^2srHz}},`$ (2)
where $`\nu _H=3.310^{15}`$ Hz, $`\alpha _\gamma 1.5`$ and the dimensionless factor $`J_{21}(z)`$ describes redshift variations of the intensity. The mean temperature, $`T_{bg}`$, of homogeneous gas can be taken as
$$T_{bg}3.5z_4^{6/7}\mathrm{\Theta }_{bg}^{4/7}10^4K,$$
$`b_{bg}=\sqrt{{\displaystyle \frac{2k_BT_{bg}}{m_H}}}24z_4^{3/7}\mathrm{\Theta }_{bg}^{2/7}\mathrm{km}/\mathrm{s},`$ (3)
$$\mathrm{\Theta }_{bg}=\frac{\mathrm{\Omega }_bh^2}{0.02}\frac{3.5}{2+\alpha _\gamma }\left(\frac{0.15}{\mathrm{\Omega }_mh^2}\right)^{3/4},z_4=\frac{1+z}{4},$$
where $`k_B\&m_H`$ are the Boltzmann constant and the mass of the hydrogen atom and $`\alpha _\gamma `$ is the power index of spectrum of the ionizing background in (2). Analyzing the observed characteristics of absorbers McDonald et al. (2001) estimate the background temperature as $`T_{bg}(2\pm 0.2)10^4K`$ at $`z2`$, what is close to (3).
At this period the gas entropy can be characterized by the function
$`F_{bg}={\displaystyle \frac{T_{bg}}{n_b^{2/3}}}=60z_4^{8/7}\mathrm{\Theta }_{bg}^{4/7}\left({\displaystyle \frac{0.02}{\mathrm{\Omega }_bh^2}}\right)^{2/3}\mathrm{keV}\mathrm{cm}^2.`$ (4)
As was shown in Demiański & Doroshkevich (2004b), the function $`F_{bg}(n_b/n_b)^{0.1}`$ only weakly depends upon variations of the expansion rate. This means that the background temperature and Doppler parameter vary as
$`T_{bg}(n_b/n_b)^{2/3},b_{bg}(n_b/n_b)^{1/3},`$ (5)
and the mildly nonlinear compression or expansion of mater occurs almost adiabatically. Thus, within compressed or slowly expanded regions with $`n_bn_b`$ we can expect that $`T_{bg}T_{bg}`$ while in rapidly expanded regions with $`n_bn_b`$ we can expect that $`T_{bg}T_{bg}`$ .
Under the assumption of ionization equilibrium of the gas,
$`{\displaystyle \frac{n_H}{n_b}}={\displaystyle \frac{\alpha _rn_b}{\mathrm{\Gamma }_\gamma }},\alpha _r4.410^{13}\left({\displaystyle \frac{10^4K}{T}}\right)^{3/4}{\displaystyle \frac{\mathrm{cm}^3}{\mathrm{s}}},`$ (6)
where $`\alpha _r(T)`$ is the recombination coefficient (Black, 1981) and $`\mathrm{\Gamma }_\gamma `$ characterizes the rate of ionization by the UV background, the fraction of neutral hydrogen is
$`x_{bg}=n_H/n_b=x_0(1+z)^{33/14},`$ (7)
$$x_0\frac{1.210^7}{\mathrm{\Gamma }_{12}(z)}\frac{\mathrm{\Omega }_bh^2}{0.02}\mathrm{\Theta }_{bg}^{3/7},\mathrm{\Gamma }_\gamma (z)=10^{12}\mathrm{\Gamma }_{12}(z)s^1,$$
$$\mathrm{\Gamma }_{12}(z)=12.6J_{21}(z)(3+\alpha _\gamma )^1.$$
The redshift variations of the rate of ionization, $`\mathrm{\Gamma }_{12}(z)`$, produced by the UV radiation of quasars were discussed by Haardt & Madau (1996) and later on tested and corrected by Demiański & Doroshkevich (2004b) with large observed sample of QSOs. For $`\alpha _\gamma 1.5`$ the expected ionization rate of hydrogen is fitted by the expression:
$`\mathrm{\Gamma }_{12}7\mathrm{exp}[(z2.35)^2/2].`$ (8)
Variations of the power index $`\alpha _\gamma `$ with time and space generate random variations of $`\mathrm{\Gamma }_{12}`$. With these parameters of background we have for the Gunn–Peterson optical depth
$`\tau _{GP}(z){\displaystyle \frac{0.34}{\mathrm{\Gamma }_{12}\mathrm{\Theta }_{bg}^{3/7}}}\left({\displaystyle \frac{1+z}{4}}\right)^{27/7}\left({\displaystyle \frac{\mathrm{\Omega }_bh^2}{0.02}}\right)^2\sqrt{{\displaystyle \frac{0.15}{\mathrm{\Omega }_mh^2}}},`$ (9)
$$\tau _{GP}(2)0.015,\tau _{GP}(5)8.2.$$
However these estimates of $`\mathrm{\Gamma }_{12}`$ should be corrected for absorption and reemission of radiation within high density clouds (Haardt & Madau 1996). This effect decreases the ionization rate of both hydrogen and helium and enhances a possible spatial variations of intensity of the UV background. According to the observational estimates of Scott et al. (2002) $`\mathrm{\Gamma }_{12}14`$ at $`z2`$ (see also Levshakov et al. 2003; Boksenberg, Sargent & Rauch 2003). On the other hand, Eq. (8) takes into account only the contribution of QSOs and, so, it underestimates the ionization rate at both $`z2`$ and $`z3`$ where contribution of other probable sources of radiation becomes important.
### 2.2 The database.
The present analysis is based on 19 high resolution spectra listed in Table 1. These spectra contain 7 770 absorbers. For further discussion we selected the sample of 7 430 absorbers with $`11.9\mathrm{lg}N_{HI}15`$ and 7 411 distances between neighboring absorbers. This sample will be used for discussion of the correlation function of initial velocity field. For detailed investigation we selected a more homogeneous sample of 6 270 absorbers and 6 251 separations with $`b5km/s`$ and we restricted errors of measurement by the conditions $`\mathrm{\Delta }\mathrm{lg}N_{HI}0.2`$ and $`\mathrm{\Delta }b0.3b`$. The chosen low limit of $`b`$ is close to the spectral resolution in this sample. To test the sample dependence of the correlation function of initial velocity field we use for comparison the sample of 14 QSOs with 4 036 absorbers investigated in Paper I and Paper II .
However, a list of absorbers depends also upon the method of line identification and, for example, for QSO 0636+680 two spectra listed in Table 1 include different number of absorbers. This example shows that the methods of line identification should be unified and improved. None the less, dispersions of absorbers characteristics discussed below are defined mainly by their broad distribution functions and by the completeness of the samples. Because of this, in this paper we discuss the scatter of only the more interesting quantitative characteristics of absorbers.
As is seen from Fig. 1, the redshift distribution of absorbers is non homogeneous and the majority of absorbers are concentrated at 2 $`z`$ 3.5 . This means that some of the discussed here characteristics of absorbers are derived mainly from this range of redshifts. Absorbers at $`z3.5`$ were identified mainly in spectra of QSOs 0000-260 and 1055+461. In this range and at $`z2`$ the statistics of lines is not sufficient.
### 2.3 Observed characteristics of absorbers
For the sample of 6 270 absorbers the redshift variations of the three mean observed characteristics of absorbers, namely, the Doppler parameter, $`b`$, the column density of neutral hydrogen, $`\mathrm{lg}N_{HI}^{}=\mathrm{lg}(N_{HI}/z_4^2)`$, and the mean separation of absorbers, $`d_{sep}^{}=d_{sep}z_4^2`$, $`z_4=(1+z)/4`$ are plotted in Fig. 1 for 1.6 $`z`$ 4. These variations are well fitted by
$$\mathrm{lg}N_{HI}^{}=\mathrm{lg}(N_{HI}/z_4^2)=13.3\pm 0.08,$$
$`b=(26\pm 2.1)\mathrm{km}/\mathrm{s},z_4=(1+z)/4,`$ (10)
$$d_v^{}=2b/H(z)z_4^{3/2}=(0.12\pm 0.01)h^1\mathrm{Mpc},$$
$$d_{sep}^{}=d_{sep}z_4^2=(1.3\pm 0.15)h^1\mathrm{Mpc},$$
respectively. For the sample of 7 411 separations we get
$$d_{sep}^{}=(1\pm 0.1)h^1\mathrm{Mpc}.$$
Detailed discussion of the observed characteristics of absorbers can be found, for example, in Kim, Cristiani & D’Odorico (2002), Kim et al. (2002, 2004).
The observed probability distribution functions, PDFs, for the Doppler parameter, $`P(b)`$, the hydrogen column density, $`P(N_{HI}/z_4^2)`$, and the absorbers separation, $`P(d_{sep}z_4^2)`$, are plotted in Fig. 2. Such choice of variables allows us to suppress the redshift evolution of the mean characteristics of absorbers. PDFs for so corrected parameters only weakly vary with redshift and, as was discussed in Paper I, for the main fraction of absorbers, these variations do not exceed 10 – 15% . These PDFs are well fitted by exponential functions
$$P_{fit}(x_{HI})0.4\mathrm{exp}(0.8x_H)+2.9\mathrm{exp}(5x_H),$$
$`P_{fit}(x_b)\{\begin{array}{cc}0.15\mathrm{exp}(2.8x_b),& bb_{rap},\\ 0.9\mathrm{exp}(2.3x_b),& bb_{rap},\end{array}`$ (13)
$$P_{fit}(x_s)3.5\mathrm{exp}(1.8x_s)\mathrm{erf}^4(\sqrt{1.8x_s})/\sqrt{x_s},$$
$$x_b=\frac{b}{b},x_{HI}=\frac{N_{HI}/z_4^2}{N_{HI}/z_4^2},x_s=\frac{d_{sep}z_4^2}{d_{sep}z_4^2}=\frac{d_{sep}^{}}{d_{sep}^{}},$$
where again $`z_4=(1+z)/4`$ and $`b_{rap}23.5km/sb_{bg}`$ (3) discriminates between absorbers situated in the increasing and decreasing parts of the PDF $`P(x_b)`$ in Fig. 2. The similarity of $`b_{rap}`$ and $`b_{bg}`$ is an independent confirmation of estimates (3). For the scatter of measured PDF, $`P(x_s)`$, around the fit (13) we have
$`P_{fit}(x_s)/P(x_s)1.04\pm 0.15,x_s2.`$ (14)
Any model of the forest has problems explaining the complex shape of the observed PDF $`P(x_b)`$ and the existence of absorbers with $`bb_{bg}`$. These absorbers are sometimes related to the unidentified metal lines what is perhaps possible for colder absorbers. However, large fraction of such absorbers ($`40\%`$) indicates that majority of them must be related to usual hydrogen clouds formed within regions with lower background temperature. Our analysis shows that statistical characteristics of these absorbers are consistent with expected ones for absorbers formed due to adiabatic compression in rapidly expanded regions (see Secs. 3.5 & 4.5). Examples of such absorbers were also found in simulations (see, e.g., Bi & Davidsen 1997; Zhang et al. 1998; Davé et al. 1999). This subpopulation can also contain some number of ”artificial” caustics (McGill 1990) .
The mean size of rapidly expanded regions, $`D_{rap}`$, can be measured as a distance along the line of sight between two absorbers with $`bb_{rap}`$ closest to a colder absorber. It increases with time as
$`D_{rap}z_4^21.6\pm 0.2h^1\mathrm{Mpc},`$ (15)
and the typical mass associated with these regions is
$`M_{rap}{\displaystyle \frac{\pi }{6}}\rho (z)\delta _mD_{rap}^310^{13}z_4^3\delta _mM_{}\left({\displaystyle \frac{\mathrm{\Omega }_m}{0.3}}{\displaystyle \frac{0.7}{h}}\right),`$ (16)
where $`\delta _m=\rho _m/\rho _m1`$ . At $`z_4`$ 1 the mass $`M_{rap}`$ is in the range of galactic masses and it rapidly increases with time. This result is consistent with the expected symmetry of positive and negative initial density perturbations what leads to formation of both galaxies and rapidly expanded regions.
Even these results allow one to obtain some inferences on the evolution of the forest:
1. Regular redshift variations of the mean observed characteristics of absorbers (10) and weak redshift dependence of PDFs (13) indicate the self similar character of absorber’s evolution over the whole range of redshifts under consideration. So, we can conclude, that at these redshifts the evolution is dominated by a balanced action of the same physical factors.
2. The complex form of the PDF $`P(b)`$ confirms that absorbers have been formed within both rapidly and moderately expanded regions. Analysis of these subpopulations of absorbers allows one to estimate some parameters of such regions. We will discuss this problem in Secs. 3.5 & 4.5 .
3. The very wide range of measured Doppler parameters, 0.2 $`b/b_{bg}`$ 5, indicates a wide variety of initial perturbations. Such a wide range of Doppler parameters is not reproduced in simulations (Meiksin et al. 2001).
4. The decline of the measured $`N_{HI}(1+z)^2`$ could be related to the retained expansion of majority of observed absorbers in the transverse directions.
5. The growth of the mean observed separation between absorbers, $`d_{sep}(1+z)^2`$ indicates the progressive decline of the number density of observed absorbers that can be related to the decrease of their hydrogen column density under the observational limit $`N_{HI}10^{12}cm^2`$.
6. Comparison of measured $`d_v(z)`$ and $`(1+z)^1d_{sep}(z)`$ (10) indicates that at $`z4`$ the expected overlapping of absorbers becomes essential and, perhaps, absence of observed absorbers in spectra of the farthest QSO (Fan et al. 2002, 2003, 2004) can be partly related to this effect.
## 3 Model of absorbers formation and evolution
### 3.1 Physical model of absorbers
Many models were proposed during the last twenty years to explain regular redshift variations of the mean observed characteristics of absorbers (10) and weak redshift dependence of PDFs (13) (see references in Rauch (1998), and in Paper I and Paper II). The simplest one connects the absorber characteristics with a suitable set of early formed equilibrium clouds. This model naturally explains the weak redshift dependence of the mean Doppler parameter, $`b`$, and the observed PDFs (13). It explains also the regular redshift variations of the mean absorber separation (10). Indeed, the mean proper free path between absorbers is
$`(1+z)^1d_{sep}(z)(1+z)^3n_{abs}(z)S_{abs}(z)^1,`$ (17)
where $`n_{abs}\&S_{abs}`$ are the number density and the surface area of such clouds orthogonal to the line of sight. In this case we have from (17)
$$n_{abs}(z)(1+z)^3,S_{abs}(z)const(z).$$
However, it is well known that the existence of such low density equilibrium clouds cannot be reasonably explained. Moreover, with this model the observed evolution of the mean column density of neutral hydrogen, $`N_{HI}(1+z)^2`$, must be related to the monotonic evolution of the UV background, $`\mathrm{\Gamma }_{12}(1+z)^2`$, what disagrees with the available estimates of $`\mathrm{\Gamma }_{12}`$.
Numerous simulations indicate that the majority of the LSS elements can be related to the extended anisotropic moderate density clouds with a complex internal structure and low density envelope. Such long lived clouds are relaxed along the shorter axis and are expanded/compressed along transverse directions. The formation of DM pancakes as an inevitable first step of evolution of small perturbations was firmly established both by theoretical considerations (Zel’dovich 1970; Shandarin & Zel’dovich 1989) and numerical simulations (Shandarin et al. 1995). The anisotropic galaxy walls such as the Great Wall are observed in the SDSS and 2dF galaxy surveys and represent examples of the Zel’dovich pancakes. This suggests that absorbers can also be linked with the Zel’dovich pancakes and the reasonable model of the forest evolution can be constructed on the basis of the Zel’dovich theory.
However, in some papers devoted to the hydrodynamical simulations of the forest evolution absorbers are often identified with unrelaxed moderate density clouds and their Doppler parameter $`b`$ is related to the gradient of velocity of infalling matter rather than to thermal velocities. Of course, observed samples include some fraction of such objects mainly among weaker absorbers. However, properties of richer absorbers with metal systems are found to be in a good agreement with a model assuming local hydrostatic equilibrium (see, e.g. Carswell, Schaye & Kim 2002; Telferet al. 2002; Simcoe, Sargent, Raugh 2002, 2004; Boksenberg, Sargent, Raugh 2003; Manning 2002, 2003 a,b; Bergeron & Herbert-Fort 2005). This means that, possibly, the fraction of simulated extended unrelaxed clouds is artificially enhanced by technical limitations (see, e.g., discussion in Meiksin, Bryan & Machacek 2001; Manning 2003 and Paper II). It is important that these simulations cannot yet reproduce well enough the characteristics of observed absorbers discussed in previous Section.
This discussion shows that an adequate physical model of the complex evolution of the forest has not been proposed yet. Indeed, the relaxation of DM pancakes leads to complex internal structure of pancakes, the adiabatic compression and/or expansion of absorbers in transverse directions is changing their overdensity and temperature, the radiative cooling and bulk heating leads to the drift of the gas entropy and overdensity but leaves unchanged the depth of the potential well formed by DM distribution. Merging of pancakes increases more strongly the depth of the potential well and the gas entropy but the overdensity of the gaseous component increases only moderately. The temperature and overdensity of the trapped gas are rearranged in accordance with the condition of hydrodynamic equilibrium across the pancake continuously. These processes imply the existence of a complicated time-dependent internal structure of absorbers.
In this paper we discuss the model of absorbers evolution based on the Zel’dvich theory. This analytical model links the self similar evolution of dominated DM component with observed characteristics of absorbers. Here we assume that:
1. The DM distribution forms a set of sheet–like clouds (Zel’dovich pancakes), their basic parameters are approximately described by the Zel’dovich theory of gravitational instability applied to the CDM or WDM initial power spectrum (DD99; DD04). The majority of DM pancakes are partly relaxed, long-lived, and their properties vary due to the successive accretion of matter, merging and expansion and/or compression in the transverse directions.
2. Gas is trapped in the gravitational potential wells formed by the DM distribution. For majority of absorbers, the observed Doppler parameter, $`b`$, traces the gas temperature and the depth of the DM potential wells. We consider the possible macroscopic motions within pancakes as subsonic and assume that they cannot essentially distort the measured Doppler parameter.
3. The gas is ionized by the UV background and for the majority of absorbers ionization equilibrium is assumed.
4. For a given temperature, the gas density within the potential wells is determined by the gas entropy created during the previous evolution. The gas entropy is changing, mainly, due to relaxation of the compressed matter, shock heating in the course of merging of pancakes, bulk heating by the UV background and local sources and due to radiative cooling. These processes slowly change the entropy and density of the trapped gas. Random variations of the intensity and spectrum of the UV background enhance random scatter of the observed properties of absorbers.
In this model we expect that the mean surface density of both DM and baryonic components of pancakes weakly varies with time and the mean density of matter compressed within absorbers decreases as $`\rho _{abs}(1+z)^\upsilon ,\upsilon 1.51`$ for colder and hotter absorbers, respectively. We cannot discriminate between evolution of the absorbers number density, $`n_{abs}`$, and their surface area, $`S_{abs}`$. However, numerical simulations demonstrate that the general tendency of evolution is a sequential growth of masses and sizes of pancakes accompanied by a drift of weaker ones under the observational limit and decrease of the comoving number density of observed absorbers. Thus, as is seen from (17), if the surface area of absorbers increases $`S_{abs}(1+z)^\kappa `$ then the number density of observed absorbers decreases as $`n_{abs}(1+z)^{3+\kappa }`$. We show that the observed evolution of mean hydrogen column density and mean absorber separations (10 & 13) coincides with theoretical expectations for Gaussian initial perturbations.
The wide range of the observed Doppler parameters, 5–10 km/s $`b`$ 100 km/s, demonstrates a complex composition of the forest. In this paper we roughly identify three subpopulations of absorbers, namely, hot absorbers formed in the course of merging and shock compression, warm absorbers formed due to adiabatic and weak shock compression in moderately expanded regions, and an unexpectedly rich subpopulation of colder absorbers. We link the measured Doppler parameter with the depth of 1D potential well formed by the compressed DM component and we neglect the contribution of macroscopic velocities what restricts the achieved precision of our approach. Non the less this model reproduces quite well the self similar evolution of absorbers and allows one to reconstruct the initial correlation function of velocities and the initial power spectrum down to very small scales.
### 3.2 The initial power spectrum and correlation functions of the initial velocity field
In this Section we summarize the main results obtained in DD99 and DD04 concerning the evolution of DM pancakes and in Sec. 3.9 we show how to improve the estimates of the correlation function of the initial velocity field, and the shape of initial power spectrum discussed in Paper II.
As a reference power spectrum of initial perturbations we take the standard CDM–like spectrum with the Harrison – Zel’dovich asymptotic,
$`P(k)={\displaystyle \frac{A^2k}{k_0^4}}T^2(\eta )D_W(\eta ),\eta ={\displaystyle \frac{k}{k_0}},k_0={\displaystyle \frac{\mathrm{\Omega }_mh^2}{\mathrm{Mpc}}},`$ (18)
where $`A`$ is the dimensionless amplitude of perturbations, $`k`$ is the comoving wave number. The transfer function, $`T(\eta )`$, and the damping factor, $`D_W(\eta )`$, describing the free streaming of DM particles were given in Bardeen et al. (1986). For WDM particles the dimensionless damping scale, $`R_f`$, and the damping factor, $`D_W`$, are
$$R_f=\frac{1}{5}\left(\frac{\mathrm{\Omega }_mh^2\mathrm{keV}}{M_{DM}}\right)^{4/3},D_W=\mathrm{exp}[\eta R_f(\eta R_f)^2],$$
where $`M_{DM}`$ is the mass of WDM particles in keV (Bardeen et al. 1986). This relation illustrates the clear dependence of characteristics of small scale perturbations on the mass of DM particles.
For the spectrum (18) the coherent lengths of velocity and density fields, $`l_v\&l_\rho `$, are expressed through the spectral moments, $`m_2\&m_0`$, (DD99) as follows:
$`l_v={\displaystyle \frac{1}{k_0\sqrt{m_2}}}={\displaystyle \frac{6.6}{\mathrm{\Omega }_mh^2}}\mathrm{Mpc}31.4h^1{\displaystyle \frac{0.21}{\mathrm{\Omega }_mh}}\mathrm{Mpc},`$ (19)
$$m_2=_0^{\mathrm{}}𝑑\eta \eta T^2(\eta )D_W(\eta )0.023,l_\rho =q_0l_v,$$
$$m_0=_0^{\mathrm{}}𝑑\eta \eta ^3T^2(\eta )D_W(\eta ),q_0=5\frac{m_2^2}{m_0}.$$
The moment $`m_0`$ depends upon the mass of dominant component of DM particles and as is shown in Sec. 4.2, $`q_010^3`$, $`m_02.4`$, $`M_{DM}`$ 1 MeV .
As was demonstrated in DD99 and DD04, the basic statistical characteristics of structure are expressed through the normalized longitudinal correlation function of the initial velocity field, $`𝐯(\stackrel{~}{𝐪})`$ ,
$`\xi _v(l_vq)=3{\displaystyle \frac{(𝐪𝐯(\stackrel{~}{𝐪}_1))(𝐪𝐯(\stackrel{~}{𝐪}_2))}{\sigma _v^2q^2}},𝐪={\displaystyle \frac{\stackrel{~}{𝐪}_1\stackrel{~}{𝐪}_2}{l_v}}.`$ (20)
Here $`\stackrel{~}{𝐪}_1\&\stackrel{~}{𝐪}_2`$ are the unperturbed coordinates of two particles at $`z=`$ 0, $`q=|𝐪|`$, and $`\sigma _v^2`$ is the velocity variance. This function is expressed through the power spectrum:
$`\xi _v(l_vq)={\displaystyle \frac{3}{m_2}}{\displaystyle _0^{\mathrm{}}}𝑑\eta \eta ^2\mathrm{cos}x{\displaystyle _\eta ^{\mathrm{}}}{\displaystyle \frac{dy}{y^2}}T^2(y)D_W(y),`$ (21)
$$\eta T^2(\eta )D_W(\eta )=\frac{\sqrt{m_2}}{3}_0^{\mathrm{}}(2\mathrm{cos}x+x\mathrm{sin}x)\xi _v(l_vq)𝑑q,$$
$$x=kl_vq,\eta =k/k_0,\xi _v(0)=1,_0^{\mathrm{}}𝑑q\xi (l_vq)=0.$$
Similar relations can be also written for any initial power spectrum $`P(k)`$ .
For the CDM – like spectrum (18) with $`q_0<10^3`$ and for the most interesting range $`0.5q`$ the velocity correlation function can be fitted as follows:
$`\xi _v(q)=\xi _{CDM}1{\displaystyle \frac{1.5q^2}{\sqrt{2.25q^4+q^2+p_0^{1.4}q^{0.6}+q_0^2}}},`$ (22)
where $`p_01.110^2`$ and $`q_0`$ was introduced in (19). Further on, we will use this function as the reference one and will compare it with observational estimates of $`\xi _v(q)`$.
As was shown in DD99 and DD04, the main DM characteristics depend upon the self similar variable
$`\zeta (q,z)={\displaystyle \frac{q^2}{4\tau ^2(z)[1\xi _v(q)]}},`$ (23)
where the ’time’ $`\tau (z)`$ describes the growth of perturbations due to the gravitational instability. For the $`\mathrm{\Lambda }`$CDM cosmological model (1), and for $`z`$ 2 we have
$`\tau (z)\tau _0\left({\displaystyle \frac{1+1.2\mathrm{\Omega }_m}{2.2\mathrm{\Omega }_m}}\right)^{1/3}{\displaystyle \frac{1}{1+z}}{\displaystyle \frac{1.27\tau _0}{1+z}},`$ (24)
and $`\tau _0`$ characterizes the amplitude of initial perturbations. It is proportional to $`\sigma _8`$, the variance of the mass within a randomly placed sphere of radius 8$`h^1`$Mpc. Latest estimates (Spergel et al. 2003, 2006; Viel et al. 2004b) for the model (1) are
$`\sigma _80.9\pm 0.1,\tau _0=0.21\sigma _8(0.19\pm 0.02){\displaystyle \frac{\sigma _8}{0.9}}.`$ (25)
### 3.3 Expected characteristics of DM absorbers
In this section we introduce (without proofs) the basic characteristics of DM pancakes as a basis for further analysis. For more details see DD99 and DD04.
#### 3.3.1 DM column density of absorbers
The fundamental characteristic of DM pancakes is the dimensional, $`\mu `$, or the dimensionless, $`q`$, Lagrangian thickness (the dimensionless DM column density) :
$`\mu {\displaystyle \frac{\rho _m(z)l_vq}{(1+z)}}={\displaystyle \frac{3H_0^2}{8\pi G}}l_v\mathrm{\Omega }_m(1+z)^2q,`$ (26)
where $`l_v`$ is the coherent length of initial velocity field (19). The Lagrangian thickness of a pancake, $`l_vq`$, is defined as the unperturbed distance at redshift $`z=0`$ between DM particles bounding the pancake (20).
As was found in DD99 and DD04 for Gaussian initial perturbations, the expected probability distribution function for the DM column density is
$`N_q(\zeta ){\displaystyle \frac{2}{\sqrt{\pi }}}e^\zeta {\displaystyle \frac{\mathrm{erf}(\sqrt{\zeta })}{\sqrt{\zeta }}},W_q(<\zeta )=\mathrm{erf}^2(\zeta ),`$ (27)
$$\zeta \frac{1}{2}+\frac{1}{\pi }0.82,q6\tau ^2\zeta \frac{2\pm 0.2}{z_4^2}10^2\left(\frac{\tau _0}{0.2}\right)^2,$$
where $`z_4=(1+z)/4`$ , $`\zeta (q,z)`$ was introduced in (23) and $`W_q(<\zeta )`$ is the cumulative probability function.
Strictly speaking, the relations (26) and (27) are valid for pancakes formed and observed at the same redshift $`z_{obs}=z_f`$ because after pancake formation the transverse expansion and compression changes its DM column density and other characteristics. However, owing to the symmetry of moderate distortions of general expansion in transverse directions, these processes do not change the statistical characteristics for majority of pancakes observed at redshift $`z_{obs}z_f`$ (DD04). This means that statistically we can consider each pancake as created at the observed redshift. However, the symmetry is distorted for rapidly expanded regions observed as absorbers with $`bb_{bg}`$ because the cross–section of strongly compressed absorbers is small and, therefore, they are rarely observed. More details are given in Secs. 3.5, 3.6 and in DD04.
#### 3.3.2 Proper sizes of absorbers
The actual thickness of a DM pancake is estimated as
$`d_{abs}={\displaystyle \frac{\mu }{\rho _m}}={\displaystyle \frac{l_vq}{(1+z)\delta _m}},\delta _m=\rho _m/\rho _m(z).`$ (28)
Here $`\delta _m`$ is the mean overdensity of compressed matter above the background density.
For the transverse size of absorbers, the expected characteristics were estimated in DD04 as follows:
$`N_{tr}(\zeta _{tr}){\displaystyle \frac{2}{\sqrt{\pi }}}\mathrm{exp}(\zeta _{tr}^2),\zeta _{tr}^2{\displaystyle \frac{q_{tr}}{6\tau ^2(z)}},`$ (29)
$$\zeta _{tr}^21/2,l_vq_{tr}3\tau ^2(z)l_v0.45z_4^2h^1\mathrm{Mpc}.$$
Here again $`z_4=(1+z)/4`$, $`z`$ is the observed redshift and $`l_vq_{tr}`$ is the expected size of absorbers at $`z=0`$.
#### 3.3.3 Absorbers separation
An important characteristic of the distribution of absorbers is their separation determined as the distance along the line of sight between centers of neighboring absorbers observed with the separation $`\mathrm{\Delta }z`$,
$`d_{sep}={\displaystyle \frac{c\mathrm{\Delta }z}{H(z)}}=5.510^3{\displaystyle \frac{\mathrm{\Delta }z}{(1+z)^{3/2}}}\sqrt{{\displaystyle \frac{0.3}{\mathrm{\Omega }_m}}}h^1\mathrm{Mpc}.`$ (30)
This separation is identical to the free path between absorbers and is quite similar to the popular descriptor
$`{\displaystyle \frac{\mathrm{\Delta }N_{abs}}{\mathrm{\Delta }z}}d_{sep}(z)^1.`$ (31)
For Gaussian initial perturbations the PDF for the absorber separations in Lagrangian space, $`d_L`$, is
$$N_L(\zeta _L)4\pi ^{1/2}\mathrm{exp}(\zeta _L)[(1+2\zeta _L)D_w(\sqrt{\zeta _L})\sqrt{\zeta _L}]$$
$`2.82\mathrm{exp}(\zeta _L)\mathrm{erf}^4(\zeta _L)/\sqrt{\zeta _L},`$ (32)
$$\zeta _L=\zeta (q_L,z),q_L=d_L/l_v,\zeta _L1.5,$$
where $`\zeta (q,z)`$ was introduced by (23) and
$$D_w(x)=_0^x𝑑y\mathrm{exp}(y^2x^2)$$
is the Dawson function.
Comparing the pancake surface density, $`\mu (z)`$, with the absorbers separation, $`d_{sep}`$, we can also estimate the fraction of matter accumulated by absorbers as
$`f_{abs}{\displaystyle \frac{\mu (1+z)}{\rho _md_{sep}}}={\displaystyle \frac{l_vq}{d_{sep}}},f_{abs}{\displaystyle \frac{q}{q_L}}0.55.`$ (33)
However, this formal estimate is of limited significance because it does not consider the complex processes of absorbers evolution. More detailed 3D analysis (DD04) shows that at a later period of absorbers evolution we can expect
$`f_{abs}0.380.44`$ (34)
The observed characteristics of absorbers are measured in the redshift space where both the proper motions of absorbers and their peculiar velocities distort the PDF (32) and, in particular, lead to the merging and artificial blending of absorbers. These processes are driven by the spatial modulations of gravitational potential formed by the large scale perturbations and can be also described in the framework of Zel’dovich approximation (DD99). However, such description depends upon the size of absorbers (28) which variations with redshift cannot be described theoretically. So, characteristics of absorbers in both real and redshift spaces can be determined only approximately.
For the PDF of absorbers separation in the redshift space we get
$`N_{sep}2.4\mathrm{exp}(1.35x_{rd})[10.85\mathrm{exp}(1.35x_{rd})]`$ (35)
$$x_{rd}=\frac{\zeta (q_{rd},z)}{\zeta (q_{rd},z)},\zeta (q_{rd},z)2.13,q_{rd}=\frac{d_{sep}}{l_v},$$
$$d_{sep}6l_v\tau ^2\zeta (q_{rd},z)1.3z_4^2h^1\mathrm{Mpc}(\tau _0/0.2)^2,$$
where $`z_4=(1+z)/4`$ and $`\zeta (q,z)\&d_{sep}`$ are defined by (23) and (30). These expressions describe correctly properties of larger separations but become unreliable for smaller separations where the influence of the proper sizes of absorbers is more important. In the real space the expression (35) also approximates the PDF of distances between neighboring absorbers but the mean value $`\zeta _{real}`$ 1.64 is smaller than that in the redshift space.
The same approach allows us to determine the expected characteristics of merged absorbers. Instead of (27), in redshift space, we get for such absorbers
$`N_{mrg}(x)1.3\mathrm{exp}(1.1x)[10.5\mathrm{exp}(2.2x)],`$ (36)
$$x=\zeta _{mrg}/\zeta _{mrg},\zeta _{mrg}=\zeta (q_{mrg},z),\zeta _{mrg}=1.22,$$
where $`\zeta (q,z)`$ was introduced in (23) and $`q_{mrg}`$ is the dimensionless DM column density of merged absorbers. This relation indicates that the population of poorer merged absorbers is suppressed and, in particular, the mean DM column density of merged absorbers is larger than that for all absorbers as given by (27).
The relations (27), (33), (35), (36) show that during the self similar period of structure evolution, when the relations (22) and (24) are valid, we can expect regular variations of the basic characteristics of absorbers such as their DM column density, $`q`$, separation, $`d_{sep}`$, and fraction of matter accumulated by absorbers, $`f_{abs}`$. These regular variations are distorted at small redshifts when the growth of perturbations is decelerated and at higher redshifts, when the blending of absorbers becomes more important.
### 3.4 Doppler parameters of absorbers
For relaxed and gravitationally confined absorbers their Doppler parameters are closely linked to the potential wells formed by the DM distribution. As is well known, for an equilibrium slab of DM the depth of its potential well is
$`\mathrm{\Delta }\mathrm{\Phi }{\displaystyle \frac{\pi G\mu ^2}{\rho (z)\delta _m}}\mathrm{\Theta }_\mathrm{\Phi }={\displaystyle \frac{3}{8}}v_0^2{\displaystyle \frac{q^2(1+z)}{\delta _m}}\mathrm{\Theta }_\mathrm{\Phi },`$ (37)
$$v_0=H_0l_v\sqrt{\mathrm{\Omega }_m}=\mathrm{1\hspace{0.17em}720}km/s\sqrt{\frac{0.15}{\mathrm{\Omega }_mh^2}}.$$
where the random factor $`\mathrm{\Theta }_\mathrm{\Phi }`$ characterizes the inhomogeneity of DM distribution across the slab and the evaporation of matter in the course of its relaxation.
Analysis of numerical simulations (Demiański et al. 2000) indicates that the relaxed distribution of DM component can be approximately described by the polytropic equation of state with the power index $`\gamma _m`$ 1.5 - 2 . Thus, for $`\gamma _m=2`$ the equilibrium density profile across the slab can be directly found and
$$\mathrm{\Theta }_\mathrm{\Phi }=4/\pi 1.3.$$
The actual distribution of DM component across a slab and the value of $`\mathrm{\Theta }_\mathrm{\Phi }`$ depends upon the relaxation process which is essentially accelerated by the process of pancake disruption into the system of high density clouds and filaments. In the course of relaxation $``$ 10 – 15% of matter is evaporated what decreases the factor $`\mathrm{\Theta }_\mathrm{\Phi }`$. This means that the expected $`\mathrm{\Theta }_\mathrm{\Phi }1`$ randomly varies from absorber to absorber.
The Doppler parameter is defined by the depth of potential well (37) and for the isentropic gas with $`p_{gas}\rho _{gas}^{5/3}`$ trapped within the well, we get
$`b^2b_{bg}^2+{\displaystyle \frac{4}{5}}\mathrm{\Delta }\mathrm{\Phi }b_{bg}^2+{\displaystyle \frac{3}{10}}v_0^2{\displaystyle \frac{q^2}{\delta _m}}(1+z)\mathrm{\Theta }_\mathrm{\Phi }.`$ (38)
Variations of the gas entropy across absorbers increase the random variations of $`\mathrm{\Delta }\mathrm{\Phi }\&b`$.
For hot absorbers with $`bb_{bg}`$ we can neglect the difference between the actual and mean background temperature and in (38) use $`b_{bg}`$ instead of $`b_{bg}`$. In such a way we link $`q^2/\delta _m`$ with $`b\&b_{bg}`$ with a reasonable precision.
### 3.5 Absorbers within rapidly expanded regions
For significant fraction of absorbers – up to 20% – the Doppler parameter is smaller then the expected mean background one (3). Such absorbers are often related to unrecognized metal lines. However, both theoretical arguments and numerical simulations show that such absorbers can also be related to hydrogen clouds formed within colder rapidly expanded regions.
The temperature of relaxed HI absorbers formed by the compression of matter (38) cannot be smaller than the background one (3). However, within low density rapidly expanded regions the background temperature given by (5) is smaller than the mean one, and in such regions the hydrogen clouds with $`b_{bg}bb_{bg}`$ can be formed. Our analysis (Sec. 4.5) shows that majority of colder absorbers could be related to such hydrogen clouds.
According to the Zel’dovich theory of gravitational instability and for Gaussian initial perturbations, regions with moderate distortions of the cosmological expansion dominate and probability to find rapidly expanded or compressed regions is exponentially small. However, the observations of galaxies at redshifts under consideration corroborate the existence of strong distortions of cosmological expansion at least on galactic scales. This means that owing to the symmetry of positive and negative initial density perturbations we should also observe rapidly expanded regions with small background density and temperature. Number of such regions exponentially decreases for larger distortions of the expansion.
Owing to the same symmetry, for majority of absorbers, fraction of absorbers with a moderate random expansion and compression in the transverse directions are close to each other and the influence of this factor weakly distorts the mean absorbers characteristics. However, within rapidly expanded regions all absorbers are adiabatically expanded in the transverse directions what, in particular, decreases the fraction of neutral hydrogen,
$$x_Hn_b/b^{3/2}\mathrm{\Gamma }_{12}n_b^{1/2}/\mathrm{\Gamma }_{12},$$
and its observed column density, and leads to a systematic drift of absorbers under the observational limit $`lgN_{HI}12`$.
Theoretical estimates show that owing to the correlation of velocity perturbations across the pancake and in the transverse directions the rate of pancakes formation within rapidly expanded regions is smaller than the mean one. For the CDM like initial power spectrum the coefficient of correlation of orthogonal velocities is $`c_v1/3`$ (DD04). In this case, for the fraction of matter, $`f_{rap}`$, and the mean column density of DM component, $`q_{rap}`$ in the rapidly expanded regions, we expect
$`f_{rap}0.33f_{abs},\zeta _{rap}=\zeta (q_{rap},z)0.5\zeta .`$ (39)
### 3.6 Characteristics of the gaseous component
The observed column density of the neutral hydrogen can be written as an integral over the line of sight through a pancake
$`N_{HI}={\displaystyle 𝑑x\rho _bx_H}=2x_H{\displaystyle \frac{n_b(z)l_vq}{1+z}}{\displaystyle \frac{0.5}{cos\theta }}.`$ (40)
Here $`x_H`$ is the mean fraction of the neutral hydrogen and $`cos\theta `$ takes into account the random orientation of absorbers and the line of sight ($`cos\theta 0.5)`$. As was noted in Sec. 3.4, we assume also that both DM and gaseous components are compressed together and, so, the column densities of baryons and DM component are proportional to each other.
Under the assumption of ionization equilibrium of the gas (6) and neglecting a possible contribution of macroscopic motions to the $`b`$-parameter ($`Tb^2`$), for the fraction of neutral hydrogen and its column density we get:
$$x_H=x_0\delta _b\beta ^{3/2}(1+z)^{33/14}\mathrm{\Theta }_x,\beta =b/b_{bg},$$
$`{\displaystyle \frac{N_{HI}}{N_0}}={\displaystyle \frac{q\delta _b}{\mathrm{\Gamma }_{12}\beta ^{3/2}}}(1+z)^{61/14},N_0=5.510^{12}cm^2\mathrm{\Theta }_H,`$ (41)
$$\mathrm{\Theta }_H=\frac{\mathrm{\Theta }_x}{\mathrm{\Theta }_{bg}^{3/7}}\frac{0.15}{\mathrm{\Omega }_mh^2}\frac{cos\theta }{cos\theta }\left(\frac{\mathrm{\Omega }_bh^2}{0.02}\right)^2,\delta _b=n_b/n_b,$$
where $`\mathrm{\Gamma }_{12}`$, $`b_{bg}`$, $`\mathrm{\Theta }_{bg}`$ and $`x_0`$ were defined in (3) and (7) and the factor $`\mathrm{\Theta }_x1`$ describes the inhomogeneous distribution of ionized hydrogen along the line of sight.
However, the overdensity of the baryonic component, $`\delta _b`$ is not identical to the overdensity of DM component, $`\delta _m`$, (see, e.g., discussion in Matarrese & Mohayaee 2002). Indeed, the gas temperature and the Doppler parameter are mainly determined by the characteristics of DM component (38) but the gas overdensity is smaller than that of DM component due to larger entropy of the gas. Moreover, the bulk heating and cooling change the density and entropy of the gas trapped within the DM potential well. These processes change the baryonic density of pancakes and we can write
$`\delta _b=\mathrm{\Theta }_b(z)\delta _m,\mathrm{\Theta }_b(z)1.`$ (42)
The factor $`\mathrm{\Theta }_b`$ should be small for absorbers formed due to adiabatic and weak shock compression because of the large difference between entropies of the background DM and the gas, and $`\mathrm{\Theta }_b`$ 1 for richer hot absorbers formed due to strong shock compression when entropies of both components are comparable.
Similarly to the proper thickness of a DM pancake (28), the thickness of a gaseous pancake is estimated as
$`d_{abs}={\displaystyle \frac{l_vq}{(1+z)\delta _b}}.`$ (43)
For adiabatically compressed absorbers it is larger then the thickness of the DM pancake but for shock compressed absorbers they are close to each other. So defined $`d_{abs}`$ can be compared with the estimated redshift thickness of absorbers determined by the observed Doppler parameter,
$`d_v=2b/H(z).`$ (44)
We can expect that, as usual, $`d_vd_{abs}`$ owing to the impact of thermal velocities. In spite of the limited precision of determination of $`d_v\&d_{abs}`$ these estimates allow us to correct the derived DM column density, $`q`$.
For long lived absorbers the influence of the bulk heating can be estimated in the same manner as it was done for characteristics of the background (Demiański & Doroshkevich 2004b). Solving the equation of thermal balance for absorbers formed at $`z=z_f`$ and observed at $`zz_f`$ we obtain for the entropy of compressed gas:
$`F_s^{3/2}(z)F_s^{3/2}(z_f)+{\displaystyle \frac{8}{7}}{\displaystyle \frac{F_{bg}^{3/2}(z)}{\beta ^{1/2}(z)}}\left[1\left({\displaystyle \frac{1+z}{1+z_f}}\right)^{3/2}\right].`$ (45)
As is seen from this relation, the bulk heating is negligible for shock compressed absorbers with $`F_s(z_f)F_{bg}(z)`$, $`\beta (z)1`$. For adiabatically compressed long lived absorbers with $`z_fz`$, $`F_s(z_f)F_{bg}(z_f)F_{bg}(z)`$ we get
$`F_s^{3/2}(z)F_{bg}^{3/2}(z)/\beta ^{1/2}(z),\delta _b\beta ^{7/2}.`$ (46)
This result demonstrates that the bulk heating is specially important for absorbers with $`\beta 1`$ formed within rapidly expanded regions. For such absorbers the merging is suppressed and their number decreases mainly owing to the drift under the observational limit.
### 3.7 Observed characteristics of absorbers
Eqs. (38) and (41) relate three independent variables, namely, $`q,\delta _m\&\delta _b`$. To find the DM column density, $`q`$, it is therefore necessary to use an additional relation which connects the basic parameters of absorbers. Here we assume that the richer absorbers with $`bb_{thr}=\beta _{thr}b_{bg}>b_{bg}`$ are formed due to shock compression and this process is accompanied by strong relaxation of compressed matter. For such absorbers $`\delta _b\delta _m`$ and their DM column density is:
$$\delta _b\delta _m\frac{v_0^2(1+z)}{b^2b_{bg}^2}q^2,bb_{thr},$$
$`q^3{\displaystyle \frac{N_{HI}\mathrm{\Gamma }_{12}}{N_0}}{\displaystyle \frac{b^2b_{bg}^2}{v_0^2}}\left({\displaystyle \frac{b}{b_{bg}}}\right)^{3/2}(1+z)^{75/14}.`$ (47)
Formation of absorbers with $`bb_{thr}=\beta _{thr}b_{bg}`$ is accompanied by adiabatic or weak shock compression of baryonic component. Assuming that the compression of baryons is described by the polytropic equation of state with $`\gamma _b=5/3`$, we can expect that for recently formed absorbers with $`bb_{thr}`$, $`F_s(z)F_{bg}(z)`$,
$`q{\displaystyle \frac{N_{HI}\mathrm{\Gamma }_{12}}{N_0}}\beta ^{3/2}(1+z)^{61/14},\delta _b\beta ^3,`$ (48)
and for long lived absorbers (46)
$`q{\displaystyle \frac{N_{HI}\mathrm{\Gamma }_{12}}{N_0}}\beta ^2(1+z)^{61/14},\delta _b\beta ^{7/2}.`$ (49)
The relations (47) – (49) determine the dimensionless column density of DM component corrected for the impact of gaseous pressure. These relations can be successfully applied to absorbers formed due to adiabatic and strong shock compression with various degrees of relaxation. However, the boundary between these limiting cases must be established a priory. So, to discriminate absorbers described by (47) and (48, 49), we use the threshold Doppler parameter, $`b_{thr}`$ which, in fact, characterizes the Mach number of the inflowing matter. Thus, absorbers with $`bb_{thr}`$,
$`b_{thr}=\beta _{thr}b_{bg}(z),\beta _{thr}1.52,`$ (50)
can be conveniently considered as the adiabatically compressed while absorbers with $`bb_{thr}`$ are considered as formed by shock compression and they contain strongly relaxed matter.
The precision of these estimates is moderate and the some uncertainties are generated by the poorly known $`\mathrm{\Gamma }_\gamma `$ and the parameters $`\mathrm{\Theta }_\mathrm{\Phi }`$ and $`\mathrm{\Theta }_H`$, which vary – randomly and systematically – from absorber to absorber. Estimates of these uncertainties can be obtained from the analysis of the derived absorbers’ characteristics. However, the main uncertainties in the estimates of $`q`$ come from the unknown $`\mathrm{cos}\theta `$, and, for strongly relaxed shock compressed absorbers, from modulations of $`b_{bg}`$.
Indeed, for adiabatically compressed absorbers $`qcos\theta `$ and Eqs.(48, 49) underestimate $`q`$ for $`\mathrm{cos}\theta \mathrm{cos}\theta =0.5`$ and overestimate it for $`\mathrm{cos}\theta \mathrm{cos}\theta `$. To reveal and to correct the most serious uncertainties we will use the condition $`d_{abs}/d_v1`$ where the real, $`d_{abs}`$, and the redshift, $`d_v`$, size of absorbers were introduced in (43, 44). For absorbers with $`d_{abs}d_v`$ we will substitute the ’true’ column density, $`q_t`$ defined by relation
$`q_t=qd_v/d_{abs},d_{abs}d_v,`$ (51)
instead of the measured one (48) or (49). Perhaps, more detailed reconstruction of the shape of absorption lines could allow us to reveal stronger deviations from the Doppler profile and, so, to identify the influence of absorbers orientation and macroscopic velocities.
For shock compressed absorbers the influence of $`\mathrm{cos}\theta `$ is not so strong as $`q\mathrm{cos}^{1/3}\theta `$ and the expected modulation of $`b_{bg}`$ is more important. Indeed, these absorbers are formed owing to the merging which is more probable in slowly expanded regions with $`b_{bg}b_{bg}`$. Thus, for absorbers with $`bb_{thr}`$, $`d_{abs}d_v`$, we will also correct $`q\&d_{abs}`$ by relation (51).
These corrections essentially improve estimates of the correlation function of initial velocity field.
### 3.8 Regular and random variations of absorbers’ characteristics
The most fundamental characteristic of absorbers is their DM column density, $`\zeta (q,z)q(1+z)^2`$ (23). It depends on the process of formation and merging of pancakes, is only weakly sensitive to the action of random factors and defines the regular redshift variations of absorbers characteristics.
For shock compressed absorbers, the evolutionary history of each pancake and the action of random factors discussed in the previous subsections are integrated in the entropy of the baryonic component,
$`S_b=\mathrm{ln}F_s(z)=S_{bg}+2/3\mathrm{ln}(\beta ^3/\delta _b).`$ (52)
If the structure of a relaxed DM pancake can be described by the polytropic equation of state with the effective power index $`\gamma _m`$ then we can introduce also the entropy of DM component, $`S_m`$. For probable value $`\gamma _m1.52`$, entropies of DM and baryonic components are quite similar to each other, $`S_mS_b`$. This means that for the strongly relaxed shock compressed absorbers the evolutionary history is characterized quite well by two functions, $`\zeta `$ and $`S_b`$.
For adiabatically compressed absorbers the baryonic entropy is identical to the background one given by (4) while the observed $`b`$ and $`N_{HI}`$ depend upon the distribution of the compressed DM component. For such absorbers the PDFs and the random scatter of observed characteristics are defined mainly by variations of the expansion rate and the background density and temperature. These characteristics as well as the entropy and overdensity of compressed DM component, $`S_m\&\delta _m`$, now cannot be derived from observational data with a reasonable reliability. This problem deserves further investigations.
### 3.9 Reconstruction of the initial power spectrum
The basic relation of Zel’dovich theory of gravitational instability can be suitably written for the difference of coordinates of two particles,
$`\mathrm{\Delta }r_i={\displaystyle \frac{l_v\tau (z)}{1+z}}[q_i/\tau (z)\mathrm{\Delta }S_i(q)/l_v],i=1,2,3,`$ (53)
where $`r_i`$ are the Euler coordinates of particles, $`l_v`$, $`\tau (z)`$ and $`q_i`$ were introduced in (19), (20) & (24) and $`S_i(q)`$ is a random displacement of a particle with respect to its unperturbed position. As is seen from this relation, for pancakes $`\mathrm{\Delta }Sl_vq/\tau (z)`$ and, so, some statistical characteristics of $`\mathrm{\Delta }S`$ can be obtained by measuring $`q/\tau (z)`$.
Using the measured redshift, $`z`$, and DM column density of absorbers, $`q`$, we determine the cumulative PDF of absorbers $`W_{obs}[>q/\tau (z)]`$ and, for each $`q/\tau (z)`$, we compute $`q`$ and $`\sigma _q^2=q^2q^2`$. For a chosen $`W_q(\zeta )`$ (27), we solve numerically the equation
$`W_{obs}[<q/\tau (z)]=W_q(\zeta )=\mathrm{erf}^2(\zeta ),`$ (54)
with respect to $`\zeta (q,\tau )`$ and, thus, we obtain the function
$`1\xi _v(q)={\displaystyle \frac{q^2}{4\tau ^2\zeta (q)}}.`$ (55)
For the most interesting range $`q/\tau (z)1,\zeta 1`$ we have
$`W_q(\zeta )4\zeta ^2/\pi ,1\xi _vq^2/\sqrt{W_q(\zeta )}.`$ (56)
The same approach can be applied to the absorbers separation, $`d_{sep}`$, and $`\zeta _{rd}=\zeta (q_{rd},z),\&q_{rd}`$ introduced by (35) . In this case instead of (54) and (56) we have
$$W_{obs}[<d_{sep}/l_v\tau (z)]=W_{sep}(<\zeta _{rd})$$
$`11.778\mathrm{exp}(0.634\zeta _{rd})[10.425\mathrm{exp}(0.634\zeta _{rd})],`$ (57)
$$1\xi _v(q_{rd})=\frac{q_{rd}^2}{4\tau ^2\zeta (q_{rd})},$$
and, for $`q_{rd}/\tau (z)1,\zeta _{rd}1`$, we get
$$W_{sep}(\zeta _{rd})0.17\zeta _{rd},1\xi _vq_{rd}^2/W_{sep}.$$
The observed functions $`P_{obs}(q)`$ and $`P_{obs}(d_{sep})`$ can be compared with corresponding expectations (27) and (35). The correlation function $`\xi _v(q)`$ derived from (55) and (57) can be compared with the reference function $`\xi _{CDM}(q)`$ (22).
## 4 Model dependent statistical characteristics of absorbers
In this Section, evolution of the basic model dependent characteristics of absorbers is discussed in the $`\mathrm{\Lambda }`$CDM cosmological model (1), and for the background temperature, $`T_{bg}`$, the Doppler parameter, $`b_{bg}`$, and the entropy function $`F_{bg}`$ given by (3) and (4) .
### 4.1 Parameters of the model
The model of absorbers discussed in Sec. 3 includes poorly known random parameters, $`\mathrm{\Theta }_{bg},\mathrm{\Theta }_\mathrm{\Phi }\&\mathrm{\Theta }_x`$, which cannot be estimated a priory, this leads to a moderate random scatter of derived absorbers characteristics. Further on we will assume that
$`\mathrm{\Theta }_{bg}=\mathrm{\Theta }_\mathrm{\Phi }=\mathrm{\Theta }_x=1.`$ (58)
However, as was discussed in Sec. 3.7, the main sources of uncertainty in (4749) are the random orientation of absorbers with respect to the line of sight, measured by $`\mathrm{cos}\theta `$ and, for strongly relaxed shock compressed absorbers, random spatial modulation of the expansion rate and the background temperature, $`b_{bg}b_{bg}`$. Large distortions of the derived parameters of absorbers can be revealed and partly corrected using the relation (51). However moderate distortions of the same parameters that can not be easily corrected restrict the precision of our approach.
The poorly known radiative ionization rate, $`\mathrm{\Gamma }_{12}`$, (8) is also an important source of uncertainty. As was noted in Sec. 2.1, now there are approximate estimates of the UV background produced by the observed QSOs but the expected ionization rate should be corrected for the absorption and reemission of UV radiation by the gas compressed within high density clouds (Haardt & Madau 1996) and for the additional emission of UV radiations by galaxies (at $`z2`$) and at $`z3`$ by poorly known sources such as Ly–$`\alpha `$ emitters (Boksenberg, Sargent & Rauch 2003; Giavalisco et al. 2004; Ouchi 2005).
Below we will describe the ionization rate by the expression
$`\mathrm{\Gamma }_{12}(z)G_0\left({\displaystyle \frac{1+z}{4}}\right)^{p_\gamma }\mathrm{exp}\left[{\displaystyle \frac{(zz_\gamma )^2}{2\sigma _\gamma ^2}}\right],`$ (59)
$$G_0=4.3,z_\gamma =1.0,\sigma _\gamma =1.58,p_\gamma =1.5,$$
where the choice of $`p_\gamma ,z_\gamma `$ and $`\sigma _\gamma `$ corrects $`\mathrm{\Gamma }_{12}(z)`$ for the impact of additional sources of radiation.
Relatively small value of $`\mathrm{\Gamma }_{12}(2)3`$ (59) is close to the observational estimates of Scott et al. (2002) and $`\mathrm{\Gamma }_{12}(5)0.2`$ is similar to the estimates of Fan et al. (2002, 2004). It shows that at $`z2`$ the UV background was probably overestimated in Haardt & Madau (1996) and Demiański & Doroshkevich (2004b) . Indeed, relatively small observed $`N_{HI}10^{13.4}cm^2`$, with $`\mathrm{lg}N_{HI}15`$, and moderate number of lines $`N_{line}=226`$ with $`lgN_{HI}15`$ in our sample, shows that the absorption of the UV background by $`HI`$ cannot be very important. However, stronger absorption by $`HeII`$ found by Levshakov et al. (2003) can distort the spectrum of the UV background and decrease $`G_0`$ and the background temperature, $`T_{bg}\&b_{bg}`$ (3). Below the fit (59) will be tested by comparing the theoretically expected and derived functions, $`\xi _v(q)`$ and $`\xi _v(q_s),\zeta (z)0.82`$ and $`l_vq(z)0.44d_{sep}(z)`$.
For such $`\mathrm{\Gamma }_{12}(z)`$ we get for the Gunn–Peterson optical depth, $`\tau _{GP}(z)`$,
$`\tau _{GP}(2)0.045f_{hom}^2,\tau _{GP}(5)5f_{hom}^2,`$ (60)
where $`f_{hom}`$ is the fraction of homogeneously distributed matter. Estimates obtained in Sec. 4.3 show that at $`z4,f_{hom}0.5`$. However, at such redshifts $`\tau _{GP}`$ is measured for regions with suppressed lines where larger $`\mathrm{\Gamma }_{12}`$ can be expected. For large $`z4`$ we can take $`f_{hom}1`$ with large scatter.
The choice of $`\mathrm{\Gamma }_{12}(z)`$ (59) coincides with
$$\tau _00.19,\sigma _80.9,$$
what agrees quite well with $`\sigma _8`$ derived by Spergel et al. (2003, 2006) and with the independent estimate (10). The final results only weakly depend upon the threshold parameter $`b_{thr}(1.52)b_{bg}`$ discriminating between adiabatically and shock compressed absorbers. Here we use $`b_{thr}1.5b_{bg}`$ .
Such choice of the model parameters allow us to obtain reasonable description of properties of absorbers for the selected sample. Variations of limits used for the sample selection lead to moderate variations of parameters (59).
### 4.2 Correlation functions of the initial velocity field
#### 4.2.1 Correlation functions derived from absorber separations
Using the method described in Sec. 3.9 and characteristics of separations between absorbers obtained in Sec. 2.2 we can estimate also the correlation function of the initial velocity field, $`\xi _v(q_{rd})`$ . This approach uses only the measured redshifts and, so, the derived function does not depend upon the measured $`b\&N_{HI}`$ and the model of absorbers discussed in Sec. 2. However, the theoretically derived PDFs (35) and (57) become unreliable at small separations and, as was noted in Sec. 2.2, already at redshifts $`z`$ 3 the blending of absorbers becomes essential what distorts the derived function $`\xi _v(q_{rd})`$.
The observed cumulative PDF, in redshift space, $`W_{obs}[q_{rd}/\tau (z)]`$, and the reconstructed correlation function $`1\xi _v(q_{rd})`$ (57) are plotted in Fig. 3 for the samples of 19 QSOs (6 251 and 7 411 separations) and 14 QSOs (3 660 separations). For these samples we have, respectively,
$$W_{obs}[<q_{rd}/\tau (z)]/W_{sep}(<\zeta _{fit})1.07\pm 0.19,1.03\pm 0.2,$$
$`{\displaystyle \frac{1\xi _v(q_{rd})}{1\xi _{fit}(q_{rd})}}0.95\pm 0.25,1\pm 0.2,`$ (61)
$$1\xi _{fit}(q_{rd})=\frac{1.5q_{rd}^2}{\sqrt{2.25q_{rd}^4+q_{rd}^2+p_s^{1.4}q_{rd}^{0.6}}},p_s=2.110^3,$$
$$0.5q1.410^3,17h^1\mathrm{Mpc}l_vq_{rd}0.03h^1\mathrm{Mpc}.$$
where $`\zeta _{fit}=\zeta (\xi _{fit},\tau )`$ is given by (23).
As is seen from Fig. 3, at $`q_{rd}10^2`$, $`l_vq_{rd}0.33h^1`$Mpc the derived correlation function coincides with the standard CDM – like one (22) for heavy DM particles, with $`q_010^3,M_{DM}`$ 1MeV . Similar results can be found also with the PDF (32) what indicates the weak sensitivity of the derived correlation function $`\zeta _{fit}`$ on the detailed shape of the used PDF.
At small scales, $`10^3q_{rd}10^2`$, we see an excess of power with respect to the reference function $`\xi _v`$ (22). This excess is caused by the deficit of observed absorbers with small separation which increases the derived function, $`1\xi _v(q_{rd})W_{sep}^1`$ (see Sec. 3.9). In this range of $`q_{rd}`$ the number of measured separations is limited, $`N_{sep}320`$ and $`N_{sep}170`$ for the samples of 19 and 14 QSOs, respectively, what decreases reliability of this result. To test impact of this factor we calculated the function $`\xi _v(q_{rd})`$ for the sample of 7 411 separations obtained without placing any restrictions on the properties of absorbers. For this extended sample the number of small separations increases up to 630 but the difference between the derived and reference functions, $`\xi _v(q_{rd})`$ and $`\xi _{CDM}(q_{rd})`$, at $`q_{rd}10^2`$ remains the same.
These results indicate that possibilities and applicability of this approach are limited. Indeed, the positions of absorbers are measured in the redshift space with a typical error $`\mathrm{\Delta }z510^5`$, $`\mathrm{\Delta }q_{rd}310^2[4/(1+z)]^{3/2}`$ comparable with the size of absorbers as measured by their Doppler parameter (10), $`(1+z)d_v/l_v310^2[4/(1+z)]^{1/2}`$ . These factors lead to the artificial blending of close absorbers, distort their characteristics and decrease reliability of the estimates (61) for $`q_{rd}10^2`$.
#### 4.2.2 Correlation functions derived from DM column density of absorbers
The same correlation function of the initial velocity field, $`\xi _v(q)`$, can be also found from estimates of the DM column density of absorbers, $`q`$. Here we use a more complex procedure of determination of $`q`$ from the observed $`z,b,\&N_{HI}`$ and poorly known $`\mathrm{\Gamma }_{12}`$ (59) what decreases its reliability. On the other hand, comparison of the functions $`\xi _v`$ found with two different approaches allows us to test the model of absorbers discussed in Secs. 3 and 4.1 .
The observed cumulative PDF, $`W_{obs}(q/\tau )`$, and the reconstructed correlation function $`1\xi _v(q)`$ are plotted in Fig. 4 for the samples of 19 QSOs with 6 270 and 7 430 absorbers and 14 QSOs (3 674 absorbers). For these samples we have, respectively,
$$W_{obs}(<q/\tau )/W_q(\zeta _{fit})1.\pm 0.07,0.9\pm 0.2,$$
$`{\displaystyle \frac{1\xi _v(q)}{1\xi _{fit}(q)}}0.95\pm 0.15,1.1\pm 0.2,`$ (62)
$$1\xi _{fit}(q)\frac{1.5q^2}{\sqrt{2.25q^4+q^2+p_q^{1.4}q^{0.6}}},p_q=0.810^3,$$
$$10^4q0.3,3h^1\mathrm{kpc}l_vq9.4h^1\mathrm{Mpc},$$
where $`\zeta _{fit}=\zeta (\xi _{fit},\tau )`$ is given by (23).
As is seen from Fig. 4, at $`q510^3`$ the derived correlation function (62) differs from the standard CDM – like one (22) . For these $`q`$ both samples of absorbers are quite representative with $`N_{abs}1590`$ and $`N_{abs}850`$ for samples of 19 and 14 QSOs, respectively. However, as was noted above, a more complex procedure of determination of $`\xi _v`$ decreases its reliability.
At these scales the difference between the derived and reference correlation functions, $`\xi _v(q)`$ and $`\xi _{CDM}`$ (22), could be mainly caused by the deficit of weaker absorbers (Sec. 3.9) because $`1\xi _v(q)W_q^{1/2}`$. As before, to test the impact of this factor we calculated the function $`\xi _v(q)`$ for the sample of 7 430 absorbers with $`\mathrm{lg}N_{HI}15`$. For such sample the difference between the derived and reference functions, $`\xi _v(q)`$ and $`\xi _{CDM}(q)`$, at $`q10^3`$ decreases and the correlation function is fitted by the expression
$`1\xi _{fit}(q)={\displaystyle \frac{1.5q^2}{\sqrt{2.25q^4+q^2+p_f^{1.4}q_{rd}^{0.6}}}},p_f=410^3.`$ (63)
The functional forms of expressions (22), (62) and (63) are identical and they differ only by the values of fit parameters $`p_q0.07p_0,p_f0.3p_0`$. This means that the difference between $`\xi _v`$ (62) and $`\xi _{CDM}`$ (22) could be mainly related to a possible incompleteness of the observed samples of weaker absorbers, to the limited precision of measurements of $`b\&N_{HI}`$ and to the limited precision of our model in describing such absorbers (see discussion in Sec. 3.7).
The correlation function (63) is quite similar to the reference one and their difference is in the range of observational errors. This fact demonstrates that probably the CDM like initial power spectrum can be traced at least down to $`q10^4`$, $`l_vq3h^1`$ kpc. This means also that $`q_010^4`$, $`m_025`$ and the effective mass of the DM particles $`M_{DM}100`$ MeV. However, the function (63) is derived from inhomogeneous samples and therefore its reliability is in question.
At larger scales, $`q510^3,l_vq0.15h^1`$ Mpc, results obtained from the analysis of both characteristics of absorbers are quite similar:
$`{\displaystyle \frac{1\xi _v(q_{rd})}{1\xi _{CDM}(q_{rd})}}1.08\pm 0.25,1.06\pm 0.17,`$ (64)
$`{\displaystyle \frac{1\xi _v(q)}{1\xi _{CDM}(q)}}1.0\pm 0.2,1.09\pm 0.16.`$ (65)
These results confirm the CDM – like type of the initial power spectrum down to scales $`l_vq0.15h^1`$Mpc what extends conclusions of Croft et al. (2002), Viel et al. (2004b); McDonald et al. (2004) and Zaroubi et al. 2005. It also demonstrates the self consistency of the adopted model of absorbers.
It is important, that similar results are found for both samples of 19 QSOs with 6 270 absorbers and of 14 QSOs with 3 674 absorbers. The samples used are compiled from spectra observed with different instruments and resolutions and the parameters of absorbers were found with different codes what increases their possible non homogeneity. Under these conditions, the stability of our results demonstrates their objectivity.
Reconstruction of the initial power spectrum with the help of the relations (21) shows that in the range of errors the measured and CDM–like power spectra are quite similar each other. At larger $`k`$, $`k/k_0100`$, there is some excess of the power but estimates become unstable because of the limited range of measured $`q`$. Investigations of this important problem should be continued with more homogeneous sample of observed spectra.
### 4.3 Statistical characteristics of the full sample of absorbers
For the sample of 6 270 absorbers the redshift variations of the mean DM column density, $`\zeta `$, the mean fraction of DM component accumulated by absorbers, $`f_{abs}`$, and the real and redshift sizes of absorbers along the line of sight, $`d_{abs}z_4^{3/2}`$, and $`d_vz_4^{3/2}`$, are plotted in Fig. 5. The PDFs for the functions $`\zeta `$ and $`d_{abs}z_4^{3/2}`$ are plotted in Fig. 6. For $``$ 1 500 absorbers of this sample the parameter $`q`$ was corrected as described in Sec. 3.7 (Eq. (51)). As was expected, the fraction of such absorbers ($``$ 25%) is close to the probability $`0\mathrm{cos}\theta 0.25`$ to find absorbers oriented along the line of sight.
The mean DM column density of absorbers, $`q`$ or $`\zeta `$, is the most stable characteristic of the sample. In principle, $`\xi `$ does not change due to the formation and merging of absorbers and due to their transverse compression and/or expansion (DD04). It depends upon the ionization rate, $`\mathrm{\Gamma }_{12}`$, the amplitude of initial perturbations, $`\tau _0`$ or $`\sigma _8`$, and upon the shape of the correlation function of initial perturbations, $`\xi _v`$, or the initial power spectrum, $`p(k)`$. It also weakly depends upon the parameter $`b_{thr}`$ used to discriminate absorbers formed by the adiabatic and shock compression. The differences between the expected and measured $`\xi `$ characterize, in fact, the completeness and representativity of the samples, the scatter of the function $`\mathrm{\Gamma }_{12}`$ and the influence of disregarded factors such as $`\mathrm{\Theta }_{bg},\mathrm{\Theta }_\mathrm{\Phi }\&\mathrm{\Theta }_x`$.
For the sample under consideration the redshift variations of the measured $`\xi (z)`$ around the mean values are moderate,
$`\zeta 0.82\pm 0.08,qz_4^2(1.8\pm 0.2)10^2.`$ (66)
At small and larger redshifts, $`z2`$ and $`z3.7`$, limited statistic of absorbers decreases reliability of our estimates. The measured $`\zeta (z)`$ is close to the theoretically expected value (27) what verifies the choice of the ionizing rate $`\mathrm{\Gamma }_{12}(z)`$ in (59) and the amplitude of initial perturbations, $`\tau _0=0.19,\sigma _80.9`$ .
The PDF of the DM column density, $`P(\zeta )`$ plotted in Fig. 6 is fitted with a scatter $``$ 10% by the function
$`P(x_\zeta )=1.2\mathrm{exp}(x_\zeta )\mathrm{erf}(\sqrt{x_\zeta })/\sqrt{x_\zeta },x_\zeta =\zeta (q,z)/\zeta ,`$ (67)
which is very close to the theoretical relation (27) . This result verifies the self consistency of the physical model used here and the assumed Gaussianity of initial perturbations.
As is well known, for the full sample the HI column density and Doppler parameter are weakly correlated, and for our sample their linear correlation coefficient is
$`R_{bHI}=[bN_{HI}^{}bN_{HI}^{}]/\sigma _b\sigma _{HI}^{}0.16,`$ (68)
where $`N_{HI}^{}=N_{HI}/z_4^2`$. At the same time, the DM column density, $`\zeta `$, is correlated with both the HI column density and Doppler parameter, and the linear correlation coefficients defined in the same manner as (68) are
$`R_{\zeta _b}0.34,R_{\zeta _HI}0.72.`$ (69)
The PDF of observed absorber separations, $`P(x_{rd})`$, plotted in Fig. 6 is fitted with a scatter $``$ 13% by the function
$`P(x_{rd})=3.56\mathrm{exp}(1.6x_{rd})\mathrm{erf}^4(\sqrt{1.6x_{rd}})/\sqrt{x_{rd}},`$ (70)
where $`x_{rd}=\zeta _{rd}(d_{sep},z)/\zeta _{rd}`$, $`\zeta _{rd}1.5`$ and $`\zeta _{rd}(d_{sep},z)`$ was introduced in (35). At small and large $`x_{rd}`$ this fit differs from the theoretically expected one for redshift space (35) but it is quite similar to the fit (32) for the PDF of separations in Lagrangian space. Perhaps, this fact can be related to the unexpectedly moderate influence of the peculiar velocities of absorbers. However, it can be partly caused by peculiarities of our samples.
Our model of absorbers (see Sec. 3) allows one to estimate roughly the real size of baryonic distribution across the absorber, $`d_{abs}`$, (43) . For the full sample we have
$`d_{abs}^{}=d_{abs}z_4^{3/2}(0.11\pm 0.01)h^1\mathrm{Mpc},`$ (71)
$$d_v^{}=d_vz_4^{3/2}(0.12\pm 0.01)h^1\mathrm{Mpc},$$
and both sizes increase with time. The PDF of the real size of baryonic pancakes is roughly fitted by:
$`P(x_d)0.6\mathrm{exp}(0.8x_d)+0.4\mathrm{exp}[(0.73x_d)^2/0.15],`$ (72)
$$x_d=d_{abs}^{}/d_{abs}^{}=d_{abs}z_4^{3/2}/d_{abs}z_4^{3/2},$$
with a scatter $`12\%`$.
For the expected mean transverse size of absorbers (29) and for model parameters (59) we have for the mean proper size
$$l_vq_{tr}(1+z)^10.45z_4^2(1+z)^1h^1\mathrm{Mpc},$$
what is consistent with recent direct estimates by Becker, Sargent & Raugh (2004) at $`z33.5`$
$$l_vq_{tr}(0.150.2)h^1\mathrm{Mpc}.$$
The exponential PDF of the transverse sizes (29) and its strong redshift dependence explain large scatter of the sizes ($`l_vq_{tr}1h^1`$ Mpc) measured in many observations of pairs of QSOs (see, e.g., discussion in Becker, Sargent & Raugh 2004).
The mean measured fraction of matter accumulated by absorbers is
$`f_{abs}(z)0.44\pm 0.07.`$ (73)
In spite of the limited applicability of the one dimensional approach (33) and the limited precision of our model of absorbers the measured fraction (73) is close to the theoretical expectation of the Zel’dovich theory (34) what verifies the choice of the model characteristics (Sec. 4.1). The weak redshift variations of this function agrees well with the self similar evolution of absorbers.
### 4.4 Adiabatically and shock compressed absorbers
To investigate the complicated evolution of absorbers in more details, we compare subpopulations of adiabatically and shock compressed absorbers. These subpopulations were separated by comparison of the measured Doppler parameter, $`b`$, with the background one $`b_{bg}`$, (3). By definition, absorbers with $`bb_{thr}=1.5b_{bg}`$ belong to the subpopulation of shock compressed and strongly relaxed absorbers, while absorbers with $`bb_{thr}=1.5b_{bg}`$ are considered as formed in a course of adiabatic or weak shock compression. This discrimination is not strict however and characteristics of subpopulations depend upon the sample used in the analysis and the parameters of the parameters of the background (359). This classification allows one to characterize reasonably well the observed absorbers and to trace their evolutionary history. For both subpopulations, redshift variations of the mean characteristics are listed in Table 2 and some of the PDFs are plotted in Fig.7.
For the sample under investigation $``$ 80% of absorbers are compressed adiabatically and they accumulate $``$ 80% of the compressed matter. These fractions, $`f_n\&f_{abs}`$, weakly vary with redshift. The PDFs, $`P(x_\zeta ),x_\zeta =\zeta /\zeta `$, plotted in Fig. 7 for both subpopulations are quite similar to each other and to the PDFs (36) and $`P(x_\zeta )`$ (67) plotted in Fig. 6 for the full sample. It demonstrates that in wide range of redshifts absorbers could be formed by both processes while the cutoff at $`\zeta =0.3\zeta `$ in the PDF for shock compressed absorbers is imposed by the method used for the discrimination of absorbers. For both samples, the correlation coefficients $`R_{\zeta _b}`$ and $`R_{\zeta _HI}`$ defined in the same manner as (68) are similar to each other. These similarities verify the generic nature of absorbers and indicate that they can be successfully combined into one sample, what is consistent with expectations of the Zel’dovich theory. For the subpopulation of shock compressed and strongly relaxed absorbers $`\zeta \zeta _{mrg}`$ (36), as expected for merged absorbers, what confirms the importance of merging in the process of absorbers’ evolution.
For both subpopulations $`b/b_{inf}1`$, what can be partly related to evaporation of high velocity particles in the course of relaxation of the compressed matter. For both subpopulations, the redshift size of absorbers along the line of sight, $`d_v`$, is larger than the real size, $`d_{abs}`$ (Table 2). The redshift dependence $`d_{abs}z^{3/2}`$ is consistent with the evolution of their redshift size. The PDFs of the real size of absorbers, $`P(x_d)`$, is plotted in Fig. 7.
For adiabatically compressed absorbers the entropy is the same as for the background while the overdensity, $`\delta _b=\beta ^3`$, is determined by $`b`$ and $`b_{bg}`$. In particular, for cold absorbers with $`bb_{bg}`$ we have $`\delta _b`$ 1.
For subpopulation of strongly relaxed and shock compressed absorbers both the entropy and the overdensity depend upon the complex evolutionary history of absorbers and cannot be described by a simple theoretical model. For such absorbers growth of the overdensity with time can be explained by the action of several factors such as successive merging and contribution of long lived absorbers formed at redshifts $`z_f`$ larger than the observed one, $`z_{obs}z_f`$ . For such absorbers the measured PDF $`P(x_\delta )`$ is plotted in Fig. 7 and it can be fitted by a superposition of two exponential functions:
$`P(x_\delta )3\mathrm{exp}(3.1x_\delta )+0.13\mathrm{exp}(0.2x_\delta ),`$ (74)
$$x_\delta =\delta _b(1+z)^{2.5}/\delta _b(1+z)^{2.5}.$$
The entropy of strongly relaxed and shock compressed baryons (52) increases with time and its PDF plotted in Fig. 7 is well fitted by the Gauss function
$`P(x_f)\mathrm{exp}[(x_f0.95)^2/0.35],`$ (75)
$$x_f=\mathrm{ln}[F_s(z)z_4^2]/\mathrm{ln}[F_s(z)z_4^2].$$
Such PDF naturally arises when the entropy is generated by the action of many random factors such as the shock waves accompanying the successive merging of absorbers.
### 4.5 Absorbers in rapidly and moderately expanded regions
As was discussed in Sec. 3.5, we expect that majority of absorbers with small Doppler parameters, $`bb_{rap}`$ 23.5 km/s, and especially with $`bb_{bg}`$ could be formed within rapidly expanded regions and, so, they can characterize some properties of these regions. Some characteristics of absorbers with $`bb_{rap}`$ are listed in Table 3 in comparison with the same characteristics of adiabatically compressed absorbers with $`b_{thr}bb_{rap}`$ situated within moderately expanded regions.
The rapidly expanded regions accumulate $`f_n`$ 50% of adiabatically compressed absorbers and $`f_{rap}`$ 30% of adiabatically compressed matter. These values practically do not depend on the redshift. The Doppler parameter, $`b`$, and the DM column density, $`\zeta _{rap}(q,z)`$, also weakly depend upon redshift. For this subpopulation we have
$`\zeta _{rap}0.5\zeta ,f_{rap}0.25f_{abs}.`$ (76)
As was noted in Sec. 2.2, the mean size of rapidly expanded regions, $`D_{rap}`$, as well as the mean separation of absorbers within these regions, $`d_{sep}`$, increase with time in the same manner as the mean separation of absorbers for the full sample (10). At $`z_41`$ the typical mass associated with rapidly expanded regions, $`M_{rap}`$ (16), is in the range of galactic masses and it increases with time $`(1+z)^3`$. This result is consistent with the expected symmetry of positive and negative initial density perturbations what leads to formation of both galaxies and rapidly expanded regions. These results are consistent with theoretically expected ones (39), what confirms the interpretation of the complex shape of PDF $`P_b(b)`$ and subpopulation of weak absorbers proposed in Sec. 3.5 .
For the subpopulation of moderately expanded absorbers variations of the Doppler parameter, $`b`$, are small, $`R_{\zeta b}1`$ and the DM column density, $`\zeta `$, depends mainly upon the hydrogen column density, $`N_{HI}`$. In contrast, for the rapidly expanded absorbers the influence of both $`b\&N_{HI}`$ are equally important. In spite of this difference, for both subpopulations the PDFs $`P(\zeta )`$ are quite similar to each other and to the PDF (67) obtained for the full sample. This fact indicates that the interaction of small and large scale perturbations changes $`\zeta `$ more strongly but only weakly influences the shape of the PDFs $`P(\zeta )`$.
For both subpopulations, the mean proper sizes of absorbers, $`d_{abs}`$, are similar (Table 3) but their PDFs are quite different. Thus, for rapidly expanded regions the PDF $`P(x_d)`$ is step–like and it is responsible for the bump at $`d_{abs}z_4^{3/2}d_{abs}z_4^{3/2}`$ in the PDF $`P(x_d)`$ plotted in Fig. 6. This distribution differs from the distribution of Doppler parameter, $`P(x_b)`$, what suggests a complex internal structure of such absorbers.
### 4.6 Absorbers and properties of the background
Some characteristics of absorbers could be used to estimate the redshift variations of the mean properties of homogeneously distributed hydrogen (see, e.g., Hui & Gnedin 1997; Schaye et al. 1999, 2000; McDonald et al. 2001). However, such estimates are inevitably approximate and their significant scatter is caused by the action of many random factors discussed above. To obtain more stable results Schaye et al. (1999, 2000) consider a cutoff at small $`b`$ in the distribution of $`b(N_{HI})`$. However, such absorbers are probably formed within rapidly expanded regions and for them both background properties and expansion rate vary randomly from absorber to absorber.
For the subpopulation of moderately expanded absorbers, we can combine Eqs. (43), (44), (38) and (48) and, in principle, connect the background temperature with the Doppler parameter of absorbers, $`b`$, and their hydrogen column density, $`N_{HI}`$. However, reliability and significance of such estimates are in question. Some restrictions on the intensity of UV background were discussed in Sec. 4.1 .
## 5 Summary and Discussion.
In this paper we continue the analysis initiated in Paper I and Paper II that is based on the statistical description of Zel’dovich pancakes (DD99, DD04). This approach allows one to connect the observed characteristics of absorbers with fundamental properties of the initial perturbations without any smoothing or filtering procedures, to reveal and to illustrate the main tendencies of structure evolution. It demonstrates also the generic origin of absorbers and the Large Scale Structure observed in the spatial distribution of galaxies at small redshifts.
We investigate the more representative sample of $``$ 6 000 absorbers what allows us to improve the physical model of absorbers introduced in Paper I and Paper II and to obtain reasonable description of physical characteristics of absorbers. The progress achieved demonstrates again the key role of the representativity of the observed samples for the construction of the physical model of absorbers and reveals a close connection between conclusions and the observational database. Further progress can be achieved with richer and more refined sample of observed absorbers.
### 5.1 Main results
Main results of our analysis can be summarized as follows:
1. For suitable parameters of the model (Sec. 4.1), the basic observed properties of absorbers and their evolution are quite successfully described by the statistical model of DM confined structure elements (Zel’dovich pancakes) with various evolutionary histories. Comparison of independent estimates of the DM characteristics of pancakes confirms the self consistency of the physical model. This model is in a good agreement with measured properties of metal systems (see, e.g. Carswell, Schaye & Kim 2002; Telferet al. 2002; Bergeron & Herbert-Fort 2005).
2. The PDFs of the DM column density and the distances between neighboring absorbers are found to be consistent with the Gaussian initial perturbations with the CDM–like initial power spectrum.
3. For the observed range of redshifts the evolution of absorbers is close to self–similar one. This implies that it leads to slow variations of mean absorber characteristics with redshift and retains their PDFs.
4. We estimate the shape of the correlation function of the initial velocity field what in turn allows us to estimate the shape of the initial power spectrum. At scales $`0.15h^1`$ Mpc both derived correlation functions, (61) and (62), reproduce the CDM–like one. This means that at such scales the power spectrum of initial perturbations is close to the standard one (18) . At smaller scales we see some differences between the derived and CDM–like correlation functions which depend upon the sample used in the analysis.
5. Analysis of variations of the Doppler parameter, $`b`$, along the line of sight demonstrates existence of rapidly expanded regions which can be considered as examples of strong negative density perturbations of galactic mass scale.
6. Our analysis shows that in the observed range of redshifts we can expect slow variations of the intensity of UV background radiation and of the ionization rate, $`\mathrm{\Gamma }_{12}`$. Our results are close to the estimates of the UV background in Haardt & Madau (1996), Scott et al. (2002), and Demiański & Doroshkevich (2004b).
### 5.2 Test of the model of absorbers
The physical model of absorbers introduced in Sec. 3 links the measured $`z`$, $`b`$ and $`N_{HI}`$ with other physical characteristics of both gaseous and DM components forming the observed absorbers. It is important that this 1D model provides us with the self consistent statistical description of the Ly–$`\alpha `$ forest although some parameters of pancakes remain unknown. Action of these parameters as well as uncertainties in the available estimates of the background temperature and the UV background radiation lead to moderate random scatter of the derived characteristics of absorbers. Fortunately, actions of these factors partly compensate each other, what allows us to obtain reasonable statistical description for majority of absorbers. The self consistency of this approach is confirmed by similarity of the functions $`\xi _v(q_{rd})`$ (61) and $`\xi _v(q)`$ (62) and by estimates of the matter fraction, $`f_{abs}`$, accumulated by absorbers (73). These functions are related to independent characteristics of absorbers obtained from measurements of their separation and their DM column density.
For richer absorbers, both the hydrostatic equilibrium of compressed matter along the shorter axis of pancakes and the close link between the gas temperature and the Doppler parameter, $`b`$, are confirmed by comparison of characteristics of the HI and metal systems (see, e.g. Carswell, Schaye & Kim 2002; Telferet al. 2002; Simcoe, Sargent & Rauch 2002, 2004; Boksenberg, Sargent and Rauch 2003; Manning 2002, 2003 a,b; Bergeron & Herbert-Fort 2005). In particlar, for 191 high resolution metal systems presented in Boksenberg, Sargent and Rough (2003) differences between the gas temperatures measured by the Doppler parameters of HI and CIV do not exceed $`25\%`$, what is comparable to the precision of measurements. Comparison of the Doppler parameters measured for HI, CIV and OVI (Carswell, Schaye & Kim 2002) verifies also their similarity and shows that as a rule the macroscopic (turbulent) velocities are subsonic. These observational results strongly support the domination of long–lived gravitationally bound and partly relaxed absorbers composed of both DM and baryonic components.
Numerical simulations show that the line width depends upon the thermal broadening, the differential Hubble flow and peculiar velocities and the relative influence of these factors varies from absorber to absorber (see, e.g., Theuns, Schaye & Haehnelt 2000; Schaye 2001). The Hubble flow is more essential for weaker absorbers and can artificially increase their Doppler parameter. For majority of absorbers in the considered 1D model the possible contribution of Hubble flow is naturally linked with the compression or expansion of pancakes in the transverse directions and depends upon the (unknown) relative orientation of absorber and the line of sight. The available observational data do not allow to discriminate between the thermal and macroscopic broadening of the lines what increases the random scatter of our results. To perform such discrimination a more detailed description of the observed line profiles is required.
One of the important problem facing the high resolution numerical simulations is the development of the methods for the more detailed reconstruction of the physical properties and revealing of links between the DM and baryonic components of observed absorbers. In particular, this includes the discrimination of the thermal and macroscopic broadening of lines, explanation of the surprisingly weak redshift dependence of the mean Doppler parameter, detection of the complex internal structure of absorbers as indicated by the observations of metal systems and so on.
However, now technical limitations restrict facilities of simulations. Thus, the small box size used eliminates the large scale part of the power spectrum and decreases the representativity of simulated sample of absorbers. As was discussed in Paper II and in Manning (2003 a,b), these factors eliminate the interaction of large and small scale perturbations and distort characteristics of the simulated absorbers. Simulations reproduce the observed transmitted flux and its main features and now they are used mainly for the surprisingly stable reconstruction of the initial power spectrum from the flux characteristics (see, e.g., Seljak et al. 2004; McDonald et al. 2004,Viel et al 2004a, b). However, the analysis of Meiksin, Bryan and Machacek (2001) shows that simulations have problems with reproduction of the observed PDFs for the column density of neutral hydrogen, $`N_HI`$, and the Doppler parameter, $`b`$, and their self similar redshift evolution.
More detailed criticism of the “Fluctuating Gunn-Petersen approximation” and the simulations of the forest can be found in Manning (2002, 2003 a,b), Paper II and references cited in these papers.
### 5.3 Properties of absorbers
Analysis of the mean absorbers characteristics performed in Sec. 4 shows that the sample of observed absorbers is composed of pancakes with various evolutionary histories. We discuss five main factors that determine evolution of absorbers after their formation. They are: the transverse expansion and compression of pancakes, the disruption of structure elements into a system of high density clouds, the merging of absorbers and the radiative heating and cooling of compressed gas. The first two factors change the overdensity of DM and gas but do not change the gas entropy. Next two factors change both the gas entropy and overdensity but do not change the DM characteristics.
The sample of observed forest can be naturally divided into subsamples of adiabatically and shock compressed absorbers formed by merging. Moreover, about half of adiabatically compressed absorbers are formed within rapidly expanded regions where the background temperature is less than mean one. So, the temperature of absorbers formed within such regions can be also less than the mean temperature of the background (3). These results illustrate the influence of some of the factors mentioned above. However, the slow variation of the mean characteristics of absorbers and their PDFs with redshift confirms that we observe the self–similar period of absorbers evolution when the action of these factors is balanced and regular variations of the UV background does not distort this balance.
For shock compressed absorbers, introduction of the DM column density, $`q`$ and $`\zeta `$, and entropy, $`S_mS_b`$, allows to discriminate between the systematic and random variations of their properties. The former ones are naturally related to the progressive growth with time of the DM column density of absorbers, $`q(z)\&\zeta (z)`$, and they can be described theoretically. On the other hand the action of random factors cannot be satisfactorily described by any theoretical model. However, in the framework of our approach, the joint action of all random factors is summarized by one random function, $`S_b`$, directly expressed through the observed parameters (52). These results alleviate the problem of description of absorbers and, perhaps, the modelling of the Ly-$`\alpha `$ forest based on the simulated DM distribution (Viel et al. 2002) .
For adiabatically compressed absorbers, the spatial distributions, entropy and overdensity of baryonic and DM components are different. Unfortunately at present these characteristics of DM component cannot be determined from observations with a reasonable accuracy. For such absorbers the baryonic entropy is identical to the background one given by (4) while the PDFs and the random scatter of observed characteristics are determined mainly by random variations of the expansion rate and the background density and temperature. For this subpopulation, the process of formation and evolution of absorbers should be investigated more thoroughly.
### 5.4 Characteristics of the initial power spectrum
The initial power spectrum of density perturbations is created at the period of inflation and its observed determination is very important for investigations of the early Universe. The amplitude and the shape of large scale initial power spectrum are approximately established by investigations of relic radiation (see. e.g, Spergel et al. 2003, 2006) and the structure of the Universe at $`z<`$ 1 detected in large redshift surveys such as the SDSS (Dodelson et al. 2002; Tegmark et al. 2004) and 2dF (Percival et al. 2001). The shape of the initial power spectrum at small scale can be tested at high redshifts where it is not so strongly distorted by nonlinear evolution (see, e.g., Croft et al. 2002; Tegmark et al. 2002; McDonald et al. 2004; Seljak et al. 2004; Zaroubi et al. 2005).
Here we retrieve the correlation function of initial velocity field, $`\xi _v`$, from direct measurements of the PDFs of fundamental characteristics of absorbers such as their separation, $`d_{sep}`$, and the DM column density, $`q`$ . Both estimates are derived in the same way and result in the same shape of the correlation function at larger scales, $`q510^3,l_vq0.15h^1`$ Mpc . At these scales the measured correlation functions coincide with the CDM–like one (22) what confirms conclusions of Croft et al. (2002), Viel et al. (2004b); McDonald et al. (2004) and Zaroubi et al. (2005) obtained at scales $`1h^1`$ Mpc.
At smaller scales the results obtained with analysis of the absorbers separation, $`d_{sep}`$, and the DM column density, $`q`$, are different and demonstrate some excess of power at scales 150 kpc $`l_vq`$ 3 kpc. Parameters of these functions and the excess depend upon the sample of absorbers and for the extended samples the derived correlation functions become quite similar to the CDM–like one. However, reliability of this result is in question due to a probable incompleteness of the extended samples.
The interpretation of these distortions is not unique because of very limited available information. As was shown in Sec. 4.2, they are sensitive to the deficit of weaker absorbers and small separations of absorbers in the sample under consideration. Therefore these distortions can be enhanced by the probable incompleteness of the observed sample created by the finite resolution of observations, blending of lines and approximate character of our analysis. If this explanation is correct than these factors restrict the presently available range of investigations to $`l_vq`$ 100 kpc . Further progress can be achieved with more refined observations of absorption spectra of QSOs and with more refined identification of absorbers in the observed spectra.
In turn, these divergences can be related to special features in the initial power spectrum at small scales. Recent WMAP measurements indicate that adiabatic Gaussian perturbations dominate on large scale (Peiris et al. 2003; Komatsu et al. 2003). However, these results do not preclude deviations from the standard CDM–like power spectrum at small scales. In particular, such deviations appear in models of the one field inflation with a complicated inflation potential (see, e.g., Ivanov, Naselsky & Novikov 1994) or multiple fields inflation (see, e.g., Polarski & Starobinsky 1995; Turok 1996). Both models generate adiabatic or isocurvature deviations from the simple CDM–like power spectrum. More detailed discussion of such models can be found, for example, in Peiris et al. (2003).
### 5.5 Absorbers as elements of the Large Scale Structure of the Universe
At redshifts $`z1.7`$ the Large Scale Structure is observed mainly as systems of absorbers in spectra of distant QSOs. Numerical simulations show that even at such redshifts we can see also high density filaments and clumps formed by “galaxies” and some of them are actually observed in spectra of QSOs as metal systems, Lyman damped and Lyman limit systems. However, available observational data do not yet allow one to characterize statistically properties of such structure elements (see, e.g., Boksenberg, Sargent & Rauch 2003).
At small redshifts, the Large Scale Structure is observed as a spatial distribution of both galaxies and neutral hydrogen. The investigation of galaxy distribution in the SDSS DR1 (Doroshkevich, Tucker, Allam & Way 2004a) results in estimates of typical parameters of galaxy walls as
$`q0.4,b320km/s,d_{sep}60h^1\mathrm{Mpc}.`$ (77)
With these data the expected column density of neutral hydrogen within the typical wall (40) is $`N_{HI}10^{11}cm^2`$ and even so spectacular object as the ’Greet Wall’ does not manifest itself through absorbers.
Our results indicate the generic link of absorbers and DM Zel’dovich pancakes and demonstrate that the embryos of walls could also be seen already at $`z3`$. Indeed, for basic parameters of subpopulation of 1 370 shock compressed absorbers with $`b30`$ km/s we have
$`d_{sep}(50\pm 11)(1+z)^2h^1\mathrm{Mpc},`$ (78)
$$q(0.4\pm 0.07)(1+z)^2,$$
what is quite similar at z=0 to that given in (77). This fact indicates that, in principle, such absorbers can be considered as embryos of wall–like elements of the Large Scale Structure of the Universe. Of course, such identification of walls observed in the galaxy distribution with elements of Ly–$`\alpha `$ forest is quite arbitrary and ignores the actual complex evolution of the LSS elements. However, it confirms generic character of the LSS evolution from richer absorbers to galaxy walls. The problem deserves further investigation first of all with more representative numerical simulations.
For the first time poor absorbers at small redshifts were observed by Morris et al. (1991, 1993) and $`1000`$ of such absorbers were found by Bahcall et al. (1993, 1996) and Jannusi et al. (1998). Some of these absorbers are identified with halos of galaxies (see, e.g., Lanzetta et al. 1995; Le Brune, Bergeron & Boisse 1996) or galaxy filaments (Penton, Shull & Stock 2002) but others are situated far from any galaxies. These observations demonstrate that the space between the LSS elements – so called ’voids’ – is not empty and contains essential fraction of baryonic and DM components of the matter.
More detailed characteristics of absorbers at $`z`$ 1 are given in Penton, Shull & Stock (2000, 2002); McLin et al. (2002) where the main absorber properties are found to be similar to those observed at high redshifts. For 79 absorbers with $`12\mathrm{lg}N_{HI}15`$, $`11km/sb80km/s`$ listed in these papers the mean absorber separation is
$`d_{sep}(10\pm 3)h^1\mathrm{Mpc}.`$ (79)
Despite the strong difference of many conditions at $`z1`$ and $`z1.5`$, these observed characteristics of absorbers are quite similar to expected ones (10) extrapolated to $`z=0`$.
### 5.6 Observed and expected evolution of the Large Scale Structure
Comparison of the expected and derived from observations characteristics of absorbers demonstrates that at the observed range of redshifts, $`1.7z4.5`$ we see the self – similar period of structure evolution. During this period the main factors determining the evolution of absorbers such as the pancake expansion, creation and merging, are balanced what leads to relatively slow evolution of the mean properties of absorbers, such as $`d_{sep}`$ and the DM column density, $`q`$ . This slow evolution is supported by slow regular variations of the UV background radiation and the ionization rate, $`\mathrm{\Gamma }_{12}`$ .
However, at small redshifts, $`z0.5`$ the growth of perturbations and merging of absorbers becomes decelerated due to the influence of the $`\mathrm{\Lambda }`$–term while expansion and disruption of absorbers remains important. This means that at such redshifts the quiet evolution of absorbers is distorted and we can expect a progressive decrease of linear density of observed absorbers with the hydrogen column density $`N_{HI}10^{12}cm^2`$ . The variations of the population of observed absorbers are also modulated by the poorly known variations of the UV background.
At larger redshifts evolution of DM pancakes is mainly driven by the shape of the initial power spectrum. For the standard CDM–like correlation function (22) we get that the self–similar evolution takes place at redshifts
$`zz_{thr}\sqrt{0.375/q_0}19\sqrt{10^3/q_0},`$ (80)
and at $`zz_{thr}`$ pancakes with $`qq_0`$ are more abundant.
However, the observational test of these expectations is quite problematic because the observed characteristics of absorbers depend also upon evolution of the background temperature and UV radiation.
### 5.7 Reheating of the Universe
Recent observations of high redshift quasars with $`z`$ 5 (Djorgovski et al. 2001; Becker et al. 2001; Pentericci et al. 2002; Fan et al. 2002, 2003, 2004) provide clear evidence in favor of the reionization of the Universe at redshifts $`z`$ 6 when the volume averaged fraction of neutral hydrogen is found to be $`f_H10^3`$ and the photo ionization rate $`\mathrm{\Gamma }_\gamma (0.020.08)10^{12}s^1`$ . These results are consistent with those expected at the end of the reionization epoch which probably takes place at $`z`$ 6.
These results can be compared with expectations of the Zel’dovich approximation (DD04). The potential of this approach is limited since it cannot describe the nonlinear stages of structure formation and, so, it cannot substitute the high resolution numerical simulations. However, it describes quite well many observed and simulated statistical characteristics of the structure such as the redshift distribution of absorbers and evolution of their DM column density. This approach does not depend on the box size, number of points and other limitations of numerical simulations (see discussion in Paper II) and it successfully augments them.
This approach shows (DD04) that at $`z`$ 6 only $``$ 3.5% of the matter is condensed within the high density clouds which can be associated with luminous objects. This value can increase up to $``$ 5 – 6% with more accurate description of the clouds collapse. The same approach also allows one to estimate the mass function of structure elements (DD04) at different redshifts. At $`z`$ 6, the mean DM mass of the clouds is expected to be $`M_{cl}10^{10}M_{}`$ and majority of clouds have masses between $`10^3M_{cl}`$ and 10 $`M_{cl}`$. The formation of low mass clouds with $`M_{cl}10^6M_{}`$ is suppressed due to strong correlation of the initial density and velocity fields at scales $`l_\rho 0.03h^1(q_0/10^3)`$ Mpc (19). However, the numerous low mass satellites of large central galaxies can be formed in the course of disruption of massive collapsed clouds at the stage of their compression into thin pancake–like objects (Doroshkevich 1980; Vishniac 1983). The minimal mass of such satellites was estimated in Barkana, Haiman & Ostriker (2001).
This means that the investigation of absorbers observed at high redshifts should be supplemented by the study of properties of dwarf isolated galaxies and discrimination between such galaxies and dwarf satellites of more massive galaxies. Such observations seem to be a perspective way to discriminate between models with one and several types of DM particles.
### Acknowledgments
This work would not have been possible without the important contribution of M. Rauch and W.L.W. Sargent who provided us with unpublished spectra of five quasars. We are deeply grateful for their permission to use their data. This paper was supported in part by the Polish State Committee for Scientific Research grant Nr. 1-P03D-014-26 and Russian Found of Fundamental Investigations grant Nr. 05-02-16302. |
warning/0506/math0506179.html | ar5iv | text | # Ideals in non-associative universal enveloping algebras of Lie triple systems
## 1. Introduction.
Given a smooth manifold $`M`$ and a point $`eM`$, a local multiplication on $`M`$ at $`e`$ is a smooth map $`U\times UM`$ where $`U`$ is some neighbourhood of $`e`$ and the point $`e`$ is a two-sided unit, that is, $`xe=ex=x`$ for all $`xU`$. If $`x`$ is sufficiently close to $`e`$, both left and right multiplications by $`x`$ are one-to-one. Therefore, there always exists a neighbourhood $`VU`$ where the operations of left and right division are defined by the identities $`a\backslash (ab)=b`$ and $`(ab)/b=a`$ respectively. Two local multiplications at the same point $`e`$ of a manifold $`M`$ are considered to be equivalent if they coincide when restricted to some neighbourhood of $`(e,e)`$ in $`M\times M`$. Equivalence classes of local multiplications are called infinitesimal loops. (Sometimes infinitesimal loops are also called local loops.)
The importance of infinitesimal loops lies in the fact that they are closely related to affine connections on manifolds. Namely, any affine connection on $`M`$ defined in some neighbourhood of $`e`$ determines a local multiplication at $`e`$. Conversely, each (not necessarily associative) local multiplication at $`e`$ defines an affine connection on some neighbourhood of $`e`$; this gives a one-to-one correspondence between germs of affine connections and infinitesimal loops. The details can be found, for example, in .
Local non-associative multiplications on manifolds can rarely be extended to global multiplications and, thus, cannot be studied directly by algebraic means. Nevertheless, any local multiplication gives rise to an algebraic structure on the tangent space at the unit element, consisting of an infinite number of multilinear operations. Such algebraic structures are known as Sabinin algebras; for associative multiplications they specialise to Lie algebras. Given a Sabinin algebra that satisfies certain convergence conditions, one can uniquely reconstruct the corresponding analytic infinitesimal loop. Therefore, Sabinin algebras may be considered as the principal algebraic tool in studying local multiplications and local affine connections.
The general theory of Sabinin algebras has so far only been developed over fields of characteristic 0. From now on we shall assume that this is the case: unless stated otherwise, all vector spaces, algebras etc will be assumed to be defined over a field $`F`$ of characteristic zero.
Many general properties of Sabinin algebras are similar to those of Lie algebras. In particular, any Sabinin algebra $`V`$ can be realised as the space of primitive elements of some ”non-associative Hopf algebra” $`U(V)`$, called the universal enveloping algebra of $`V`$. The operations in $`V`$ are naturally recovered from the product in $`U(V)`$. Just as in the Lie algebra case, the universal enveloping algebras of Sabinin algebras have Poincaré–Birkhoff–Witt bases. If a Sabinin algebra $`V`$ happens to be a Lie algebra, $`U(V)`$ is precisely the usual universal enveloping algebra of the Lie algebra $`V`$.
The definition of a Sabinin algebra involves an infinite number of multilinear operations that satisfy rather complicated identities; we refer to , or for the precise form of these. However, additional conditions imposed on a local multiplication may greatly simplify the structure of the corresponding Sabinin algebra. For example, the associativity condition implies that only one of all the multilinear operations is non-zero; the identities of a Sabinin algebra specialise to the identities defining a Lie algebra, that is, antisymmetry and the Jacobi identity. If a local multiplication satisfies the Moufang law
(1)
$$a(b(ac))=((ab)a)c\text{and}((ca)b)a=c(a(ba)),$$
the corresponding Sabinin algebra is a Malcev algebra. A vector space with a bilinear skew-symmetric operation (bracket) is called a Malcev algebra if the bracket satisfies
$$[J(a,b,c),a]=J(a,b,[a,c])$$
where $`J(a,b,c)=[[a,b],c]+[[b,c],a]+[[c,a],b]`$ denotes the jacobian of $`a,b`$ and $`c`$.
Imposing the left Bol identity
$$a(b(ac))=(a(ba))c,$$
on the local multiplication, we obtain the structure of a left Bol algebra on the tangent space to the unit. A left Bol algebra is a vector space with one bilinear and one trilinear operation, denoted by $`[,]`$ and $`[,,]`$ respectively. The ternary bracket must satisfy the following relations:
$`[a,a,b]=0`$
$`[a,b,c]+[b,c,a]+[c,a,b]=0`$
$`[x,y,[a,b,c]]=[[x,y,a],b,c]+[a,[x,y,b],c]+[a,b,[x,y,c]].`$
The binary bracket is required to be skew-symmetric and should satisfy
$$[a,b,[x,y]]=[[a,b,x],y]+[x,[a,b,y]]+[x,y,[a,b]]+[[a,b],[x,y]].$$
Bol algebras generalise Malcev algebras. Indeed, in any Malcev algebra a ternary bracket can be defined by
$$[a,b,c]=[[a,b],c]\frac{1}{3}J(a,b,c).$$
With this additional operation a Malcev algebra becomes a Bol algebra.
Another important subclass of Bol algebras are Lie triple systems; these are the Bol algebras whose binary bracket is identically equal to zero. Lie triple systems arise as tangent spaces to smooth local Bruck loops (also known as K-loops). These loops, in addition to the left Bol identity, satisfy the identity
$$(ab)^1=a^1b^1.$$
where $`x^1`$ is shorthand for $`e/x`$; see . Lie triple systems play a prominent role in the theory of symmetric spaces since a symmetric space can be given the structure of a local Bruck loop at any point.
Identities satisfied in an infinitesimal loop can be translated into identities satisfied in the universal enveloping algebra of the corresponding Sabinin algebra. In particular, the universal enveloping algebra $`U(M)`$ of a Malcev algebra $`M`$ is a non-associative bialgebra that satisfies the linearisations
$$a_{(1)}(y(a_{(2)}z))=((a_{(1)}y)a_{(2)})z$$
and
$$((ya_{(1)})z)a_{(2)}=y(a_{(1)}(za_{(2)}))$$
of (1). Here we use Sweedler’s notation for the comultiplication: $`\mathrm{\Delta }(a)=a_{(1)}a_{(2)}`$. Since $`M`$ coincides with the subspace of all primitive elements of $`U(M)`$ we have $`\mathrm{\Delta }(a)=a1+1a`$ for any $`aM`$, and, hence, $`a(yz)+y(az)=(ay)z+(ya)z`$ and $`(ya)z+(yz)a=y(az)+y(za)`$, or, equivalently
(2)
$$(a,y,z)=(y,a,z)=(y,z,a).$$
Therefore, $`M`$ lies in the generalised alternative nucleus $`\mathrm{N}_{\mathrm{alt}}(U(M))`$ of $`U(M)`$. (The subset $`\mathrm{N}_{\mathrm{alt}}(A)`$ of an algebra $`A`$ consists of all $`aA`$ that satisfy (2) for any $`y,zA`$). The product on $`M`$ is recovered as $`[a,b]=abba`$ in $`U(M)`$.
The universal enveloping algebra $`U(V)`$ of a Bol algebra $`V`$ satisfies the identity
(3)
$$a_{(1)}(y(a_{(2)}z))=(a_{(1)}(ya_{(2)}))z.$$
Since $`V`$ coincides with the primitive elements of $`U(V)`$, for any $`aV`$ and $`y,zU(V)`$ we have that
(4)
$$(a,y,z)=(y,a,z).$$
This is equivalent to saying that $`V`$ is contained in the left generalised alternative nucleus $`\mathrm{LN}_{\mathrm{alt}}(U(V))`$ of the algebra $`U(V)`$. The binary and the ternary products on $`V`$ are recovered by
$$[a,b]=abba\text{and}[a,b,c]=a(bc)b(ac)c(ab)+c(ba)$$
in $`U(V)`$. It is known that for any algebra $`A`$
$$\mathrm{LN}_{\mathrm{alt}}(A)=\{aA|(a,x,y)=(x,a,y)x,yA\}$$
is a Lie triple system with $`[a,b,c]=a(bc)b(ac)c(ab)+c(ba)`$.
Universal enveloping algebras for Malcev, Bol and general Sabinin algebras have been introduced only recently; their properties are still waiting to be explored. It might be tempting to assume that the theory of universal enveloping algebras for Lie algebras can be extended rather painlessly to the case of general Sabinin algebras, especially since many aspects of the theory are known to generalise well. However, it turns out that some very basic properties, such as the abundance of ideals in the universal enveloping algebras of Lie algebras, fail to hold in the general non-associative case. In particular, we shall see that while the properties of Malcev and Bol algebras, discussed above, may look similar, this similarity does not extend too far.
The motivation for this paper is the following version of Ado’s Theorem for Malcev algebras that appeared in :
###### Theorem 1.
For any finite–dimensional Malcev algebra $`M`$ over a field of characteristic $`2,3`$ there exists a unital finite–dimensional algebra $`A`$ and a monomorphism of Malcev algebras $`\iota :M\mathrm{N}_{\mathrm{alt}}(A)`$.
One is prompted to ask whether a similar statement holds for other classes of Bol algebras, in particular, for Lie triple systems. Given a finite dimensional Lie triple system $`V`$, one could ask whether it is contained as a subsystem of $`\mathrm{LN}_{\mathrm{alt}}(A)`$, with $`ab=ba`$ for all $`a,bV`$, for some finite dimensional unital algebra $`A`$. It is easy to see that this happens if and only if there exists an ideal of finite codimension in $`U(V)`$ which intersects $`V`$ trivially. Our answer shows that for Lie triple systems the situation is very different from the case of Lie or Malcev algebras:
###### Theorem 2.
Let $`A`$ be a finite dimensional unital algebra over a field $`F`$ of characteristic 0 and $`V`$ — a Lie triple system contained as a subsystem in $`\mathrm{LN}_{\mathrm{alt}}(A)`$ such that $`ab=ba`$ for all $`a,bV`$. Assume that $`A`$ is generated by $`V`$ as a unital algebra. Then $`V`$ is nilpotent and $`A`$ decomposes (as a vector space) into a direct sum of a nilpotent ideal and a central subalgebra without nonzero nilpotent elements.
Note that we do not claim that embeddings mentioned in Theorem 2 do exist for all nilpotent Lie triple systems.
Examples suggest that the ideals in the universal enveloping algebras of Lie triple systems are even scarcer than it is implied by Theorem 2.
###### Conjecture 3.
The only proper ideal of the universal enveloping algebra of a simple Lie triple system is its augmentation ideal.
We shall verify the above conjecture in several cases by direct calculations in Poincaré–Birkhoff–Witt bases.
For each Lie triple system $`V`$ there exists a canonically defined Lie algebra $`_S(V)`$, called the Lie envelope of $`V`$ of which $`V`$ is a subsystem. The Poincaré–Birkhoff–Witt Theorem allows to identify the algebra $`U(V)`$ with a subspace of $`U(_S(V))`$. Motivated by analogy with Bruck loops, we shall show how the multiplication on $`U(_S(V))`$ can be modified to become compatible with the non-associative multiplication on $`U(V)`$.
The paper is organised as follows. The next section is auxiliary; it is a loose collection of various properties of Bol algebras and Lie triple systems. Section 3 contains the proof of Theorem 2. The construction of the universal enveloping algebra of a Lie triple system via its Lie envelope is given in Section 4. Finally, in Section 5 we present some evidence for Conjecture 3.
We have made no attempt to make this paper self-contained. We refer to for the properties of the universal enveloping algebras of Malcev algebras, to — for Bol algebras and to — for general Sabinin algebras. The paper of Lister is the general reference for Lie triple systems; the questions of nilpotency are treated in .
About the notation: we shall often write ”L.t.s.” for ”Lie triple system”. As usual, the true meaning of ”non-associative” is ”not necessarily associative”; however ”non-nilpotent” stands for ”not nilpotent”. The notations $`L_x`$ and $`R_x`$ are used to denote the multiplication by $`x`$ on the left and on the right respectively; the sum $`L_a+R_a`$ is denoted by $`T_a`$. The product $`a(a(\mathrm{}(aa))`$ will be written simply as $`a^n`$. The left, middle and right associative nuclei of an algebra $`A`$ are denoted by $`N_l(A)`$, $`N_m(A)`$ and $`N_r(A)`$ respectively, while $`\mathrm{Z}(A)`$ is the notation for the center of $`A`$. (Recall that the left associative nucleus of $`A`$ is the set of all $`aA`$ such that $`(a,y,z)=0`$ for arbitrary $`y,zA`$; the right and the middle associative nuclei are defined similarly.) By $`algX`$ (or $`alg_1X`$) we denote the subalgebra (unital subalgebra, respectively) generated by the subset $`XA`$.
## 2. Some properties of the enveloping algebras for Bol algebras and Lie triple systems.
###### Lemma 4.
Let $`(V,[,,],[,])`$ be a Bol algebra. For $`a,bV`$ such that $`[a,b]=0`$, the map $`[L_a,L_b]`$ is a derivation of $`U(V)`$.
Recall that a ternary derivation of an algebra $`A`$ is a triple $`(d_1,d_2,d_3)`$ of linear maps such that
$$d_1(xy)=d_2(x)y+xd_3(y)$$
for all $`x,yA`$. The set $`Tder(A)`$ of all ternary derivations of $`A`$ is a Lie algebra with the obvious bracket. It is clear that if $`d_1(1)=d_2(1)=d_3(1)=0`$ then $`d_1=d_2=d_3`$ is a derivation of $`A`$.
###### Proof of Lemma 4.
Notice that the identity (4) can be written as $`(L_a,T_a,L_a)Tder(U(V))`$ and, as a consequence,
$$([L_a,L_b],[T_a,T_b],[L_a,L_b])Tder(U(V)).$$
Evaluating both commutators at $`1`$, we observe that $`[L_a,L_b](1)=[a,b]=0=[T_a,T_b](1)`$, so $`[L_a,L_b]=[T_a,T_b]`$ is a derivation of $`U(V)`$. ∎
###### Lemma 5.
Let $`(V,[,,],[,])`$ be a Bol algebra. For $`a,bV`$ such that $`[a,b]=0`$ and any $`xU(V)`$
$$[L_a,L_b](x)=2(a,b,x).$$
###### Proof.
The identity (4) with $`y=b`$ gives $`L_aL_b+L_bL_a=L_{ab+ba}=2L_{ab}`$. Therefore, $`[L_a,L_b](x)=L_aL_b(x)(2L_{ab}L_aL_b)(x)=2(a,b,x)`$. ∎
###### Lemma 6.
Let $`(V,[,,],[,])`$ be a Bol algebra. For any $`aV`$
$$L_{a^n}L_{a^m}=L_{a^{n+m}}.$$
###### Proof.
See Proposition 38 in . ∎
For any Sabinin algebra $`V`$, the universal enveloping algebra is an H-bialgebra. That is, $`U(V)`$ is a non-associative unital bialgebra equipped with two bilinear maps, $`\backslash :U(V)\times U(V)U(V)`$ and $`/:U(V)\times U(V)U(V)`$ such that
$`{\displaystyle x_{(1)}\backslash (x_{(2)}y)}`$ $`=ϵ(x)y=`$ $`{\displaystyle x_{(1)}(x_{(2)}\backslash y)\text{ and}}`$
$`{\displaystyle (yx_{(1)})/x_{(2)}}`$ $`=ϵ(x)y=`$ $`{\displaystyle (y/x_{(1)})x_{(2)}}.`$
The behaviour of these maps with respect to the comultiplication $`\mathrm{\Delta }`$ and the counit $`ϵ`$ is expressed by
$$\mathrm{\Delta }(x\backslash y)=x_{(1)}\backslash y_{(1)}x_{(2)}\backslash y_{(2)},\mathrm{\Delta }(y/x)=y_{(1)}/x_{(1)}y_{(2)}/x_{(2)}$$
and
$$ϵ(x\backslash y)=ϵ(x)ϵ(y),ϵ(y/x)=ϵ(x)ϵ(y).$$
Fix an ordered basis $`\{a_i\}_{i\mathrm{\Lambda }}`$ of $`V`$, with $`\mathrm{\Lambda }`$ being the index set. The algebra $`U(V)`$ then has the Poincaré–Birkhoff–Witt basis
$$\{a_{i_1}(a_{i_2}(\mathrm{}(a_{i_{n1}}a_{i_n})\mathrm{}))|i_1\mathrm{}i_n\text{ and }n\}.$$
The algebra $`U(V)`$ is filtered by $`U(V)=_nU(V)_n`$ with
$$U(V)_n=spana_1(a_2(\mathrm{}(a_{m1}a_m))|a_1,\mathrm{},a_mV,mn.$$
The degree of an element of $`U(V)`$ with respect to this filtration is defined in the obvious way. The corresponding graded algebra $`GrU(V)`$ is isomorphic to $`Sym(V)`$, the symmetric algebra on $`V`$.
Let $`(V,[,,])`$ be a Lie triple system, and $`U(V)`$ — its universal enveloping algebra. The automorphism $`aa`$ of $`V`$ extends to an automorphism $`S:U(V)U(V)`$.
###### Lemma 7.
Let $`(V,[,,])`$ be a L.t.s. and $`U(V)`$ — its universal enveloping algebra. Then for any $`aV`$ we have that $`[a,U(V)_n]U(V)_{n1}`$.
###### Proof.
Let $`x=a_1(a_2(\mathrm{}(a_{n1}a_n)\mathrm{}))U(V)_n`$ with $`a_1,\mathrm{},a_nV`$. Since $`GrU(V)`$ is isomorphic to $`Sym(V)`$, $`[a,x]`$ belongs to $`U(V)_n`$. On the other hand, $`S([a,x])=[a,(1)^nx]=(1)^{n1}[a,x]`$. Therefore, $`[a,x]U(V)_{n1}`$. ∎
The automorphism $`S`$ notably simplifies the left division $`\backslash `$ on $`U(V)`$.
###### Proposition 8.
Let $`(V,[,,])`$ be a L.t.s. For all $`x,yU(V)`$
$$x\backslash y=S(x)y\text{and}S(x)=x\backslash 1=1/x.$$
###### Proof.
Let us prove that $`S(x_{(1)})x_{(2)}=ϵ(x)1`$. To this end we observe that this is a linear relation, so we only have to verify it on a set of elements spanning the vector space $`U(V)`$, for instance, $`\{1\}\{a^n|aV\}`$ with $`a^n=a(\mathrm{}(aa))`$. We have $`S(a_{}^{n}{}_{(1)}{}^{})a_{}^{n}{}_{(2)}{}^{}=_{k=0}^n\left(\genfrac{}{}{0pt}{}{n}{k}\right)S(a^k)a^{nk}=_{k=0}^n\left(\genfrac{}{}{0pt}{}{n}{k}\right)(1)^ka^n=0=ϵ(a^n)`$, as desired.
From (3) and $`S(x_{(1)})x_{(2)}=ϵ(x)1`$ we obtain
$$x_{(1)}(S(x_{(2)})(x_{(3)}y))=(x_{(1)}(S(x_{(2)})x_{(3)}))y=(x_{(1)}ϵ(x_{(2)}))y=xy.$$
By the definition of $`\backslash `$ we have
$$S(x_{(1)})(x_{(2)}y)=x_{(1)}\backslash (x_{(2)}(S(x_{(3)})(x_{(4)}y)))=x_{(1)}\backslash (x_{(2)}y)=ϵ(x)y$$
so
$$S(x)y=S(x_{(1)})(x_{(2)}(x_{(3)}\backslash y))=ϵ(x_{(1)})x_{(2)}\backslash y=x\backslash y.$$
With $`y=1`$ we get $`S(x)=x\backslash 1`$, and from $`S(x_{(1)})x_{(2)}=ϵ(x)1`$ we also get $`S(x)=(S(x_{(1)})x_{(2)})/x_{(3)}=ϵ(x_{(1)})1/x_{(2)}=1/x`$. ∎
Proposition 8 ensures that $`U(V)`$ satisfies the linearisation of the equations defining a Bruck loop. Therefore, the linearisation of any identity satisfied by Bruck loops will hold in $`U(V)`$. Consider, for instance, the so-called precession map $`\delta _{a,b}:c(ab)\backslash (a(bc))`$. For a Bruck loop this map is known to be an automorphism . Linearising this result we obtain
###### Corollary 9.
Let $`(V,[,,])`$ be a L.t.s. The map $`\delta _{x,y}:U(V)U(V)`$ given by
$$\delta _{x,y}(z)=(x_{(1)}y_{(1)})\backslash (x_{(2)}(y_{(2)}z))$$
satisfies
$$\delta _{x,y}(wz)=\delta _{x_{(1)},y_{(1)}}(w)\delta _{x_{(2)},y_{(2)}}(z).$$
The maps $`\delta _{x,y}`$ reflect the lack of associativity in $`U(V)`$. They satisfy
(5)
$$(x_{(1)}y_{(1)})\delta _{x_{(2)},y_{(2)}}(z)=x(yz).$$
Clearly, $`\mathrm{\Delta }(\delta _{x,y}(z))=\delta _{x_{(1)},y_{(1)}}(z_{(1)})\delta _{x_{(2)},y_{(2)}}(z_{(2)})`$. Thus,
(6)
$$\delta _{x,y}(V)V$$
and in general
$$\delta _{x,y}(U(V)_n)U(V)_n.$$
The maps $`\delta _{x,a}`$ and $`\delta _{a,x}`$ are derivations of $`U(V)`$ for any $`aV`$. In fact, $`\delta _{a,b}(x)=(a,b,x)`$ and $`\delta _{a,b}(c)=\frac{1}{2}[a,b,c]`$ for any $`a,b,cV`$ and $`xU(V)`$.
The following statement is a direct analogue of the corresponding result for Bruck loops .
###### Proposition 10.
Let $`(V,[,,])`$ be a L.t.s. Then the left and the middle associative nuclei of $`U(V)`$ coincide:
$$N_l(U(V))=N_m(U(V)).$$
###### Proof.
The identity (3) implies
(7)
$$x_{(1)}\left((S(x_{(2)})y)(x_{(3)}z)\right)=\left(x_{(1)}((S(x_{(2)})y)x_{(3)})\right)z.$$
If $`y`$ is in $`N_m(U(V))`$, the left-hand side of (7) is equal to
$$x_{(1)}\left(S(x_{(2)})(y(x_{(3)}z))\right)=y(xz).$$
On the other hand, the right-hand side of (7) can be re-written as
$$\left(x_{(1)}(S(x_{(2)})(yx_{(3)}))\right)z=(yx)z$$
and, hence, $`y(xz)=(yx)z`$ for all $`x,zU(V)`$. Therefore, $`N_m(U(V))N_l(U(V))`$.
Similarly, notice that (3) also implies
$$x_{(1)}\left((yS(x_{(2)}))(x_{(3)}z)\right)=\left(x_{(1)}((yS(x_{(2)}))x_{(3)})\right)z.$$
For $`yN_l(U(V))`$ one concludes that $`x(yz)=(xy)z`$ for all $`x,zU(V)`$ and, hence, that $`N_l(U(V))N_m(U(V))`$. ∎
###### Lemma 11.
Let $`(V,[,,])`$ be a L.t.s. and $`A`$ — a quotient of $`U(V)`$. If $`aV`$ satisfies $`[L_a,L_b]=0`$ for all $`bV`$, then $`a\mathrm{Z}(A)`$.
###### Proof.
For any $`xA`$ we have $`L_xalg_1L_b|bV`$. This can be established by induction on the degree of $`x`$ with respect to the PBW filtration that $`A`$ inherits from $`U(V)`$, using the fact that $`L_{by+yb}=L_bL_y+L_yL_b`$ for any $`yA`$ and $`bV`$.
Since $`[L_a,L_b]=0`$ for all $`bB`$ we have that $`[L_a,L_x]=0`$ for any $`xA`$, so $`a(xy)=x(ay)`$ for any $`x,yA`$. Setting $`y=1`$ we get that $`ax=xa`$ for any $`xA`$. Therefore, $`(xy)a=a(xy)=x(ay)=x(ya)`$ and $`aN_r(A)`$. This can also be expressed by saying that the triple $`(R_a,0,R_a)`$ belongs to $`Tder(A)`$.
The identity (4) implies that $`(L_a,T_a,L_a)`$ is also in $`Tder(A)`$. Since $`R_a=L_a`$, it follows that $`(2L_a,2L_a,0)Tder(A)`$ and thus $`aN_l(A)`$. Similarly, $`(0,2R_a,2L_a)Tder(A)`$ implies that $`aN_m(A)`$ and, therefore, $`a\mathrm{Z}(A)`$. ∎
## 3. Nonexistence of ideals of finite codimension
In this section $`(V,[,,])`$ will be a Lie triple system and $`U(V)`$ — the non-associative universal enveloping algebra of $`V`$. For any $`a,b,cV`$ we have
$$[L_a,L_b](c)=a(bc)b(ac)=[a,b,c]\text{and}[a,b]=0$$
in $`U(V)`$. The map $`[L_a,L_b]`$ is a derivation of $`U(V)`$ and $`L_{ax+xa}=L_aL_x+L_xL_a`$ for any $`aV,xU(V)`$ by (4).
Let $`A`$ be a finite-dimensional unital algebra and $`\mathrm{LN}_{\mathrm{alt}}(A)`$ — its left generalized alternative nucleus. We are interested in the existence of monomorphisms of L.t.s.
(8)
$$\iota :V\mathrm{LN}_{\mathrm{alt}}(A)$$
such that $`\iota (a)\iota (b)=\iota (b)\iota (a)`$ for any $`a,bV`$. By the universal property of $`U(V)`$ such a map induces a homomorphism $`\phi :U(V)A`$. The kernel of $`\phi `$ is an ideal of finite codimension whose intersection with $`V`$ is trivial.
Let $`S_2`$ be the two-dimensional simple L.t.s. generated by $`e,f`$ with
(9)
$$[e,f,e]=2e\text{ and }[e,f,f]=2f.$$
###### Lemma 12.
With $`e,f`$ as above,
$$[e^n,f]=n(n1)e^{n1}.$$
holds in $`U(S_2)`$.
###### Proof.
Observe that $`fe^n=f(ee^{n1})=e(fe^{n1})[L_e,L_f](e^{n1})=e(fe^{n1})2(n1)e^{n1}`$. Repeating with $`fe^{n1}`$ we obtain
$`fe^n`$ $`=`$ $`e^nf2((n1)+(n2)+\mathrm{}+1)e^{n1}`$
$`=`$ $`e^nfn(n1)e^{n1},`$
Any semisimple L.t.s. contains a copy of $`S_2`$, see . (This may be compared to the fact that any semisimple Lie algebra contains a copy of $`sl_2`$.)
###### Proposition 13.
If $`(V,[,,])`$ is a semisimple L.t.s., then the only proper ideal of $`U(V)`$ that has finite codimension is the augmentation ideal $`\mathrm{ker}ϵ`$.
###### Proof.
Given a proper ideal $`I`$ of $`U(V)`$ whose codimension is finite, the set $`V_0=IV`$ is an ideal of the L.t.s. $`V`$. Therefore, there exists another ideal $`V_1`$ with $`V=V_0V_1`$ (see ). Both $`V_0`$ and $`V_1`$ are semisimple L.t.s., so either $`V_1=0`$, or there exists a subsystem $`spane,fV_1`$ with multiplication as in (9). In the first case we have that $`\mathrm{ker}ϵ`$, the ideal generated by $`V`$, is contained inside $`I`$ and, hence, since the codimension of $`\mathrm{ker}ϵ`$ is 1, they are equal.
Assume now that we are in the second case. Since any finite–codimensional proper ideal $`I`$ of $`U(V)`$ contains an element of the form $`p(e)=\alpha _01+\alpha _1e+\mathrm{}+\alpha _{n1}e^{n1}+e^n`$ with $`n>1`$, then, by Lemma 12, it also contains $`[[p(e),f],f],\mathrm{}],f]=n!(n1)!e`$. Therefore, $`eI`$ which, by definition of $`V_1`$, is not possible. ∎
Proposition 13 shows that embeddings of the type (8) do not exist for semisimple L.t.s. Since any L.t.s. decomposes (as a vector space) as the direct sum of a semisimple subsystem and a solvable ideal (see ), it is clear that such embedding might only exist for solvable L.t.s. We shall prove that, in fact, $`V`$ must be nilpotent.
Let us denote the map $`c[a,b,c]`$ by $`D_{a,b}`$. The vector space $`_S(V)=spanD_{a,b}|a,bVV`$ is a Lie algebra (see ) with the bracket
(10)
$$[a,b]=D_{a,b}\text{and}[D_{a,b},c]=[a,b,c].$$
This Lie algebra is called sometimes the Lie envelope of $`V`$. It is $`_2`$–graded with even part $`_S^+(V)=spanD_{a,b}|a,bV`$ and odd part $`_S^{}(V)=V`$.
Given any unital algebra $`A`$ generated, as a unital algebra, by a subsystem $`V`$ of $`\mathrm{LN}_{\mathrm{alt}}(A)`$ with $`[a,b]=0`$ for any $`a,bV`$, we shall often consider the Lie algebra $`(V)`$ generated by $`\{L_a|aV\}`$. Usually, no explicit mention of $`A`$ will be needed. Since $`[L_a,L_b]`$ is a derivation of $`A`$ (see the proof of Lemma 4) and $`A`$ is generated by $`V`$, it follows that $`(V)=span[L_a,L_b]|a,bVspanL_a|aV`$. The algebra $`(V)`$ is isomorphic to $`_S(V)`$ by $`aL_a`$ and $`D_{a,b}[L_a,L_b]`$.
It is a simple exercise to check that over algebraically closed fields of characteristic zero, the only solvable non-nilpotent two-dimensional L.t.s. is $`R_2=FaFb`$ with
(11)
$$[a,b,a]=b\text{and}[a,b,b]=0.$$
###### Lemma 14.
Let $`V`$ be a solvable non-nilpotent L.t.s. Then there exists a homomorphic image of $`V`$ which contains a subsystem isomorphic to $`R_2`$.
###### Proof.
For $`V`$ a solvable L.t.s, $`_S(V)`$ is a solvable Lie algebra (see ). The solvability of $`_S(V)`$ implies that there exists a non-zero $`v_S(V)`$ and a homomorphism of Lie algebras $`\lambda :_S(V)F`$ such that
(12)
$$[x,v]=\lambda (x)v$$
for any $`x_S(V)`$. Observe that $`_S^+(V)[_S(V),_S(V)]`$ and hence $`\lambda (_S^+(V))=0`$.
Write $`v`$ as a sum of its even and odd components: $`v=D+b`$ with $`D_S^+(V)`$ and $`b_S^{}(V)`$. The odd part of the identity (12) with $`x_S^+(V)`$ implies that $`[V,V,b]=0`$. Setting $`x=aV`$ in (12) gives
$$D_{a,b}=\lambda (a)D$$
as the even part, and
$$D(a)=\lambda (a)b$$
as the odd part.
Assume that $`\lambda `$ is not identically equal to zero. Then we can choose $`aV`$ with $`\lambda (a)=1`$. For such $`a`$ we have that $`D=D_{a,b}`$ and $`D(a)=b`$ so $`[a,b,a]=b`$. Since $`[V,V,b]=0`$, the subspace $`spana,b`$ is a subsystem of $`V`$ isomorphic to $`R_2`$.
Now, if $`\lambda `$ happens to be identically equal to zero, it follows that $`D=0`$ and $`b0`$ (since $`v`$ is non-zero), and that $`[V,b,V]=0`$. Hence, the one-dimensional subspace $`spanb`$ is contained in the centre of $`V`$. The L.t.s. $`V/spanb`$ is solvable non-nilpotent (see ) and its dimension is lower than the dimension of $`V`$. The result in this case can be obtained by induction. ∎
###### Proposition 15.
Given a non-nilpotent L.t.s. $`V`$ and an ideal $`I`$ of finite codimension in $`U(V)`$, the intersection $`IV`$ is non-zero.
###### Proof.
Without loss of generality we may assume that $`V`$ is solvable. By Lemma 14 there exists an ideal $`V_0`$ and elements $`a,bV`$ such that $`V_0spana,b`$ is a subsystem of $`V`$ with $`[a,b,a]bmodV_0`$ and $`[a,b,b]0modV_0`$.
By (5), in $`U(V)`$ we have $`x(yz)=x_{(1)}y_{(1)}\delta _{x_{(2)},y_{(2)}}(z)`$. With $`x=a^n`$, $`y=cV`$ and $`z=a`$ we obtain
$$a^n(ca)=\{\begin{array}{c}a^{n+1}c=aa^nc\hfill \\ \\ a^nca+(a^n)_{(1)}\delta _{(a^n)_{(2)},c}(a)\hfill \\ a^nca+na^{n1}\delta _{a,c}(a)modU(V)_{n1}\hfill \end{array}$$
where the last congruence follows from (6). Hence $`[a^nc,a]na^{n1}\delta _{a,c}(a)\frac{n}{2}a^{n1}[a,c,a]modU(V)_{n1}`$. After $`n`$ commutations we get
$$[\mathrm{}[[a^nc,a],a],\mathrm{},a]=(1)^n\frac{n!}{2^n}[a,[a,[\mathrm{},[a,c,a],\mathrm{}],a],a]$$
where we have replaced the congruence modulo $`U(V)_0=F`$ by the equality because both sides lie inside $`\mathrm{ker}ϵ`$. In the particular case $`c=b`$ we have $`[\mathrm{}[[a^nb,a],a],\mathrm{},a]=\frac{n!}{2^n}(b+v_0)`$ with $`v_0V_0`$.
Any finite-codimensional ideal $`I`$ contains an element of the form $`p(a)=\alpha _01+\alpha _1a+\mathrm{}+\alpha _{n1}a^{n1}+a^n`$. It also contains $`p(a)b`$ and $`[\mathrm{}[p(a)b,a],\mathrm{},a]`$ where the commutator is taken $`n`$ times. Therefore, $`I`$ also contains the nonzero element $`\frac{n!}{2^n}(b+v_0)`$. ∎
We have seen that faithful representations of the type (8) can only exist for nilpotent L.t.s. It turns out that for nilpotent L.t.s. these representations, if exist, have very specific structure. Name, assuming that in (8) the algebra $`A`$ is generated by $`\iota (V)`$, we shall prove that there exists a nilpotent ideal $`R`$ such that $`A/R`$ is a commutative associative algebra over $`F`$ with no nontrivial nilpotent elements. First, we need some lemmas.
###### Lemma 16.
Let $`A`$ be a finite-dimensional unital algebra, $`a\mathrm{LN}_{\mathrm{alt}}(A)`$ and $`L_a=(L_a)_s+(L_a)_n`$ — the Jordan–Chevalley decomposition of $`L_a`$ in $`End(A)`$. Then there exist $`a_s,a_n\mathrm{LN}_{\mathrm{alt}}(A)`$, the semisimple and nilpotent parts of $`a`$, with $`(L_a)_s=L_{a_s}`$ and $`(L_a)_n=L_{a_n}`$.
###### Proof.
Recall that given $`(d,d^{},d^{\prime \prime })Tder(A)`$, its semisimple and nilpotent parts can be calculated componentwise: $`(d,d^{},d^{\prime \prime })_s=(d_s,d_s^{},d_s^{\prime \prime })`$ and $`(d,d^{},d^{\prime \prime })_n=(d_n,d_n^{},d_n^{\prime \prime })`$, where both $`(d_s,d_s^{},d_s^{\prime \prime })`$ and $`(d_n,d_n^{},d_n^{\prime \prime })`$ are also ternary derivations. Recall also that $`(d,d^{},d)Tder(A)`$ if and only if $`d=L_a`$ and $`d^{}=T_a`$ with $`a\mathrm{LN}_{\mathrm{alt}}(A)`$. Now, for any $`a\mathrm{LN}_{\mathrm{alt}}(A)`$ we have that $`((L_a)_s,(T_a)_s,(L_a)_s)`$ and $`((L_a)_n,(T_a)_n,(L_a)_n)Tder(A)`$, which implies that $`(L_a)_s=L_{a_s}`$ and $`(L_a)_n=L_{a_n}`$ for some $`a_s,a_n\mathrm{LN}_{\mathrm{alt}}(A)`$. ∎
Let us complete $`V`$ inside $`A`$ by adding the semisimple and nilpotent parts of all its elements; it turns out that such completion retains some fundamental properties of $`V`$:
###### Lemma 17.
Let $`A`$ be a finite dimensional unital algebra. Given any subsystem $`V\mathrm{LN}_{\mathrm{alt}}(A)`$ such that
* $`V`$ generates $`A`$ as a unital algebra,
* $`[a,b]=0`$ for all $`a,bV`$,
* $`V`$ is nilpotent,
there exists in $`\mathrm{LN}_{\mathrm{alt}}(A)`$ a subsystem $`\widehat{V}`$ containing $`V`$ and satisfying i), ii) and iii), and such that $`a_s,a_n\widehat{V}`$ for any $`a\widehat{V}`$. Moreover, $`a_s\mathrm{Z}(A)`$ for any $`a\widehat{V}`$ and $`\{a_n|a\widehat{V}\}`$ is an ideal of $`\widehat{V}`$.
###### Proof.
Since $`V`$ generates $`A`$ and $`[a,b]=0`$ for any $`a,bV`$, the Lie algebra $`(V)`$, generated by $`\{L_a|aV\}`$ is isomorphic to $`_S(V)`$. By the latter algebra, and hence the former, is nilpotent.
By the properties of the Jordan–Chevalley decomposition (see ) $`(ad_{L_a})_s=ad_{L_{a_s}}`$ and $`(ad_{L_a})_n=ad_{L_{a_n}}`$. The operators $`ad_{L_{a_s}}`$ and $`ad_{L_{a_n}}`$ can be expressed as polynomials in $`ad_{L_a}`$ with zero constant term. In particular, $`ad_{L_{a_s}}`$ leaves $`(V)`$ stable with a nilpotent action. By the semisimplicity of $`ad_{L_{a_s}}`$ this means that $`[L_{a_s},(V)]=0`$. Hence $`a_s\mathrm{Z}(A)`$ by Lemma 11.
As $`(V)`$ is nilpotent, there exists a basis of $`A`$ where $`(V)`$ is represented by upper triangular matrices. Hence, for any $`a,bV`$ the operator $`L_{a_s+b_s}`$ is semisimple, while $`L_{a_n+b_n}`$ is nilpotent. Moreover, $`a_s+b_s\mathrm{Z}(A)`$ implies that $`[L_{a_s+b_s},L_{a_n+b_n}]=0`$. By the uniqueness of the Jordan–Chevalley decomposition we obtain that $`(L_{a+b})_s=L_{a_s}+L_{b_s}`$ and $`(L_{a+b})_n=L_{a_n}+L_{b_n}`$. In particular, $`(a+b)_s=a_s+b_s`$ and $`(a+b)_n=a_n+b_n`$.
Let $`\widehat{V}=\{a_s+b_n|a,bV\}`$. By the previous, $`\widehat{V}`$ is a vector subspace of $`\mathrm{LN}_{\mathrm{alt}}(A)`$ and, since $`(a_s+b_n)_s=a_s`$ and $`(a_s+b_n)_n=b_n`$, $`\widehat{V}`$ contains the semisimple and nilpotent components of its elements. We also know that $`a_s\mathrm{Z}(A)`$ for any $`a\widehat{V}`$.
Given $`a,a^{},a^{\prime \prime },b,b^{},b^{\prime \prime }V`$ we have that
$$[a_s+b_n,a_s^{}+b_n^{}]=[b_n,b_n^{}]=[b_s+b_n,b_s^{}+b_n^{}]=[b,b^{}]=0,$$
so $`\widehat{V}`$ satisfies ii). Moreover,
$`[a_s+b_n,a_s^{}+b_n^{},a_s^{\prime \prime }+b_n^{\prime \prime }]`$ $`=`$ $`[L_{a_s+b_n},L_{a_s^{}+b_n^{}}](a_s^{\prime \prime }+b_n^{\prime \prime })`$
$`=`$ $`[L_{b_n},L_{b_n^{}}](a_s^{\prime \prime }+b_n^{\prime \prime })`$
$`=`$ $`[b_n,b_n^{},b_n^{\prime \prime }]`$
$`=`$ $`[b,b^{},b^{\prime \prime }]`$
implies that $`\widehat{V}`$ is a subsystem of $`\mathrm{LN}_{\mathrm{alt}}(A)`$ and that $`[\widehat{V},\widehat{V},\widehat{V}][V,V,V]`$. In terms of the lower central series for $`\widehat{V}`$ and $`V`$ (see ) this says that $`\widehat{V}^1V^1`$. Assuming that $`\widehat{V}^nV^n`$, we have $`\widehat{V}^{n+1}=[\widehat{V}^n,\widehat{V},\widehat{V}]+[\widehat{V},\widehat{V},\widehat{V}^n][V^n,\widehat{V},\widehat{V}]+[\widehat{V},\widehat{V},V^n][V^n,V,V]+[V,V,V^n]V^{n+1}`$. The nilpotency of $`\widehat{V}`$ follows from this observation and the nilpotency of $`V`$.
Finally, the left multiplication operator by $`[a,b,c]`$ is obtained as the commutator $`[[L_a,L_b],L_c]`$; in an adequate basis of $`A`$ it is represented as a commutator of upper triangular matrices. Therefore, it is nilpotent and $`[a,b,c]=[a,b,c]_n`$. Since $`[\widehat{V},\widehat{V},\widehat{V}][V,V,V]`$ it follows that $`[\widehat{V},\widehat{V},\widehat{V}]\{a_n|a\widehat{V}\}`$. In particular, the latter set is an ideal of $`\widehat{V}`$. ∎
###### Lemma 18.
Let $`A`$ be a finite-dimensional unital algebra and let $`V`$ be a subsystem of $`\mathrm{LN}_{\mathrm{alt}}(A)`$. Assume that
* $`a=a_n`$ for any $`aV`$,
* $`[a,b]=0`$ for all $`a,bV`$.
Then the subalgebra generated by $`V`$ is nilpotent.
###### Proof.
Assume, as before, that $`A`$ is generated by $`V`$ as a unital algebra.
There exists an element of $`V`$ that lies in the centre of $`A`$. Indeed, the nilpotency of $`V`$ implies that $`(V)`$ consists of nilpotent transformations , which, in turn, implies that the centre of $`(V)`$ is non-zero. Given $`0D+L_a\mathrm{Z}((V))`$ with $`D^+(V)`$, for any $`bV`$ the equality $`0=[D+L_a,L_b]=L_{D(b)}+[L_a,L_b]`$ implies that $`D=0`$ and $`[L_a,L_b]=0`$. Therefore $`0a\mathrm{Z}(A)`$ by Lemma 11.
We shall use induction on the dimension of $`V`$. The case $`dimV=0`$ is obvious. Given $`V`$ with $`dimV=n+1`$, choose $`0aZ(A)V`$ as above and consider the ideal $`aA`$. The quotient algebra $`A/aA`$ is generated, as a unital algebra, by the quotient $`(V+aA)/aA`$ of $`V`$. Thus we can apply the hypothesis of induction to conclude that $`algV+aA/aA=algV/aA`$ is nilpotent.
Let us denote the ideal $`algV`$ by $`A_0`$, and the linear span of all products of $`N`$ elements of $`A_0`$, regardless of the order of the parentheses, by $`A_0^N`$. From the nilpotency of $`algV/aA`$ we deduce that there exists $`N`$ such that $`A_0^NaA`$. Moreover, any product involving $`2N`$ elements of $`A_0`$ lies in the ideal $`aA_0`$, since is of the form $`u_1u_2`$ where at least one of the factors involves at least $`N`$ elements and, therefore, lies in $`A_0^NaA`$, and the other factor belongs to $`A_0`$.
Let us fix $`N`$ such that $`A_0^NaA_0`$ and prove by induction that $`A_0^{N^k}a^kA_0`$. Since $`aV`$ is nilpotent, this will imply that $`A_0`$ is nilpotent, as desired. Assume that $`A_0^{N^{k1}}a^{k1}A_0`$. Any product of $`N^k`$ elements in $`A_0`$ can be written as a product of $`N`$ factors, each belonging to $`A_0`$, and at least one of them lying in $`A_0^{N^{k1}}a^{k1}A_0`$. Since $`a^{k1}`$ is in the centre of $`A`$, the whole product lies in $`a^{k1}A_0^Na^kA_0`$. ∎
Finally, we are in the position to prove Theorem 2.
###### Proof of Theorem 2.
By Lemma 17 we can assume that $`V`$ contains the semisimple and nilpotent components of all its elements. Let $`Q=alga_s|aV\mathrm{Z}(A)`$ and let $`R`$ be the ideal generated by $`\{a_n|aV\}`$. Clearly $`A=Q+R`$.
For any nilpotent element $`xQ`$, $`L_x`$ belongs to $`alg_1L_{a_s}|aV`$. This algebra is abelian and all its elements are semisimple transformations. But $`x\mathrm{Z}(A)`$ implies that $`L_x`$ is nilpotent so $`L_x=0`$ and $`x`$ must be zero. Hence $`Q`$ is a commutative associative finite dimensional algebra without nonzero nilpotent elements.
Since $`a_s\mathrm{Z}(A)`$, it follows that $`A=Qalg_1a_n|aV`$. We can apply Lemma 18 to the algebra $`alg_1a_n|aV`$ and the subsystem $`\{a_n|aV\}`$ to conclude that $`alga_n|aV`$ is nilpotent. The ideal $`R`$ decomposes as $`R=Qalga_n|aV`$, so it is also nilpotent. Its nilpotency implies that $`QR=0`$. ∎
## 4. The universal enveloping algebras of a L.t.s. and its Lie envelope
The following construction is based on the known construction of a Bruck loop starting from a group whose every element has a square root. Namely, any such group with the product $`gh=g^{\frac{1}{2}}hg^{\frac{1}{2}}`$ becomes a Bruck loop. Observe that the linearisation of the identity $`g=r(g)r(g)`$ with $`r(g)=g^{\frac{1}{2}}`$ in an H–bialgebra reads as $`x=r(x_{(1)})r(x_{(2)})`$ for some map $`r`$.
Let $`L`$ be a Lie algebra over a field $`F`$ of characteristic $`2`$.
###### Lemma 19.
The linear map $`q:U(L)U(L)`$ defined by $`xx_{(1)}x_{(2)}`$ is bijective.
###### Proof.
Consider the Poincaré–Birkhoff–Witt filtration $`U(L)=_{n0}U_n`$ of $`U(L)`$. Given $`a_1,\mathrm{},a_nL`$,
$$q(a_1\mathrm{}a_n)2^na_1\mathrm{}a_nmodU_{n1}.$$
Since $`q`$ preserves the filtration, it follows that it is bijective on each $`U_n`$. ∎
Let $`r`$ be the inverse of $`q`$. Clearly, for any $`xU(L)`$ we have that $`x=r(x)_{(1)}r(x)_{(2)}`$. Furthermore, $`q`$ being a coalgebra isomorphims implies that $`r`$ is also a coalgebra isomorphism. Therefore,
$$x=r(x_{(1)})r(x_{(2)})$$
The product on $`U(L)`$ can be modified with the help of the map $`r`$ as follows:
$$xy=r(x_{(1)})yr(x_{(2)}).$$
With this product $`U(L)`$ becomes a unital non-associative algebra. In fact, since $`r`$ is a homomorphism of coalgebras, $`U(L)`$ carries the structure of an H–bialgebra.
###### Lemma 20.
For all $`x,y`$ in $`U(L)`$
$$r(x_{(1)}(yx_{(2)}))=r(x_{(1)})r(y)r(x_{(2)}).$$
###### Proof.
Indeed,
$`{\displaystyle x_{(1)}(yx_{(2)})}`$ $`=`$ $`{\displaystyle r(x_{(1)})r(y_{(1)})x_{(3)}r(y_{(2)})r(x_{(2)})}`$
$`=`$ $`{\displaystyle r(x_{(1)})r(y_{(1)})r(x_{(2)})r(x_{(3)})r(y_{(2)})r(x_{(4)})}`$
$`=`$ $`{\displaystyle \left(r(x_{(1)})r(y)r(x_{(2)})\right)_{(1)}\left(r(x_{(1)})r(y)r(x_{(2)})\right)_{(2)}}`$
which proves the lemma. ∎
###### Proposition 21.
The algebra $`(U(L),)`$ satisfies
* $`x_{(1)}(y(x_{(2)}z))=(x_{(1)}(yx_{(2)}))z`$.
* $`ab=ba`$ for any $`a,bL`$.
* $`a(bc)b(ac)=\frac{1}{4}[[a,b],c]`$ for any $`a,b,cL`$.
###### Proof.
We shall only check part i); it follows from Lemma 20 by
$$\begin{array}{c}x_{(1)}(y(x_{(2)}z))=r(x_{(1)})r(y_{(1)})r(x_{(2)})zr(x_{(3)})r(y_{(2)})r(x_{(4)})\hfill \\ \hfill =r(x_{(1)}(yx_{(2)}))_{(1)}zr(x_{(1)}(yx_{(2)}))_{(2)}=(x_{(1)}(yx_{(2)}))z.\end{array}$$
Given a L.t.s. with the product $`[,,]`$ and a scalar $`\mu `$, the new product $`[,,]^{}=\mu ^2[,,]`$ also defines a L.t.s. that is isomorphic to the original L.t.s. under $`x\mu x`$.
###### Corollary 22.
Let $`V`$ be a L.t.s. and $`_S(V)`$ — the Lie envelope of $`V`$. The unital subalgebra of $`(U(_S(V)),)`$ generated by $`V`$ is isomorphic to the universal enveloping algebra of $`V`$ considered as a Bol algebra with the trivial binary product.
###### Proof.
Define $`[a,b,c]^{}=\frac{1}{4}[a,b,c]`$ and let $`Q`$ the subalgebra of $`(U(L),)`$ generated by $`V`$. The universal property of $`U(V,[,,]^{})`$ together with Proposition 21 implies that there exists an epimorphism from $`U(V,[,,]^{})`$ to $`Q`$. Since $`a_1(\mathrm{}(a_{n1}a_n))a_1\mathrm{}a_nmodU_{n1}`$ with $`a_1,\mathrm{},a_nV`$, it follows that $`Q`$ admits a PBW–type basis. The epimorphism from $`U(V,[,,]^{})`$ to $`Q`$ maps the PBW basis of $`U(V,[,,]^{})`$ to this basis, so it is an isomorphism. However, as $`(V,[,,]^{})`$ and $`(V,[,,])`$ are isomorphic, their universal enveloping algebras also are. ∎
## 5. Ideals in the enveloping algebras of simple L.t.s.
###### Lemma 23.
Assume that $`V`$ is a simple L.t.s. satisfying the following condition: all elements of the universal enveloping algebra $`U(V)`$ that commute with $`V`$, are of the form $`c+x`$ where $`c`$ is a scalar and $`x`$ is in $`V`$. Then the only proper ideal of $`U(V)`$ is the augmentation ideal.
###### Proof.
Suppose the conditions of the lemma are satisfied. Let $`IU(V)`$ be an ideal, and take some $`rI`$. There exists an element $`xV`$ such that $`r^{}=rxxr0`$. It is clear that $`r^{}I`$ and $`\mathrm{deg}r^{}<\mathrm{deg}r`$, where the degree is taken with respect to the PBW filtration. Hence, $`I`$ necessarily contains a nonzero element $`u`$ of degree at most 1. If $`u`$ is a scalar, then $`I=U(V)`$. If $`\mathrm{deg}u=1`$, the space of all linear combinations of (possibly iterated) brackets containing $`u`$, is an ideal of $`V`$ and, hence, coincides with $`V`$. All these brackets are in $`I`$, therefore, $`I`$ contains $`V`$. ∎
If $`a`$ is an element of $`V`$ and $`rU(V)_n`$, the commutator $`arra`$ belongs to $`U(V)_{n1}`$. In fact, it is possible to write an explicit formula for the terms of degree $`n1`$ in this commutator.
###### Lemma 24.
Let $`\{x_k\}`$ be a basis for $`V`$ and $`rU(V)`$ — a monomial in the $`x_k`$. Then
(13)
$$arra=\frac{1}{2}\underset{i,j}{}[a,x_i,x_j]\frac{}{x_i}\frac{}{x_j}r+\text{lower degree terms}.$$
Here the partial derivative $`/x_i`$ of a non-associative monomial is defined by setting $`/x_i(uv)=u/x_i(v)+/x_i(u)v`$ with $`/x_i(x_j)=1`$ if $`i=j`$ and 0 otherwise.
###### Proof.
The vector space $`U(V)_n/U(V)_{n1}`$ is spanned by classes of elements of the form $`b^n`$ with $`bV`$, so it is sufficient to verify (13) for $`p=b^n`$.
Modulo terms of degree $`n2`$ and smaller we have
$`ab^nb^na`$ $`=a(bb^{n1})b^na`$
$`=[L_a,L_b](b^{n1})+b(ab^{n1})b^na`$
$`={\displaystyle \underset{i+j=n2}{}}b^i([L_a,L_b](b)b^j)+b(ab^{n1})b^na`$
$`=(n1)[L_a,L_b](b)b^{n2}+b(ab^{n1}b^{n1}a)`$
$`=(n1)[L_a,L_b](b)b^{n2}+(n2)[L_a,L_b](b)b^{n2}+\mathrm{}+[L_a,L_b](b)b^{n2}`$
$`={\displaystyle \frac{n(n1)}{2}}[L_a,L_b](b)b^{n2}.`$
The last expression coincides with the right-hand side of (13). ∎
Let $`x,y,z`$ be a set of generators for the Lie algebra $`so(3)`$ with $`[x,y]=z`$, $`[y,z]=x`$ and $`[z,x]=y`$. We shall consider $`so(3)`$ as a simple L.t.s. by setting $`[a,b,c]=[[a,b],c]`$. Let $`\stackrel{~}{S}_2`$ be the 2-dimensional subsystem spanned by $`x`$ and $`y`$. Over the complex numbers $`\stackrel{~}{S}_2`$ is isomorphic to the L.t.s. $`S_2`$ mentioned in Section 2; the isomorphism is given by $`e=x+y\sqrt{1}`$, $`f=x+y\sqrt{1}`$.
###### Proposition 25.
Both $`so(3)`$ and $`\stackrel{~}{S}_2`$ satisfy Conjecture 3.
###### Proof.
The products of the form $`z^n(x^py^q)`$ with $`n,p,q`$ non-negative integers, form a basis for the universal enveloping algebra of $`so(3)`$ considered as a Lie triple system. In our case, (13) reads as
$$\begin{array}{c}z^n(x^py^q)zzz^n(x^py^q)\hfill \\ \hfill =\frac{n(p+q)}{2}z^{n1}(x^py^q)+\frac{p(p1)}{2}z^{n+1}(x^{p2}y^q)+\frac{q(q1)}{2}z^{n+1}(x^py^{q2})+\mathrm{}\end{array}$$
where the omitted terms are of degree $`n+p+q2`$ and smaller. Similarly,
$$\begin{array}{c}z^n(x^py^q)xxz^n(x^py^q)\hfill \\ \hfill =\frac{n(n1)}{2}z^{n2}(x^{p+1}y^q)\frac{p(n+q)}{2}z^n(x^{p1}y^q)+\frac{q(q1)}{2}z^n(x^{p+1}y^{q2})+\mathrm{}\end{array}$$
and
$$\begin{array}{c}z^n(x^py^q)yyz^n(x^py^q)\hfill \\ \hfill =\frac{n(n1)}{2}z^{n2}(x^py^{q+1})+\frac{p(p1)}{2}z^n(x^{p2}y^{q+1})\frac{q(n+p)}{2}z^n(x^py^{q1})+\mathrm{}\end{array}$$
Now, suppose that there exists an element $`r`$ of the universal enveloping algebra of $`so(3)`$ considered as a Lie triple system, of degree $`N>1`$, which commutes with $`x,y`$ and $`z`$. This element has the form
$$r=\underset{n+p+q=N}{}\alpha _{n,p,q}z^n(x^py^q)+\text{lower degree terms.}$$
The requirement that $`rzzr`$ has no terms of degree $`N1`$ imposes linear conditions on the coefficients $`\alpha _{n,p,q}`$, similar conditions come from $`rxxr`$ and $`ryyr`$. Explicitly, these conditions are as follows:
$`(p+q)(n+2)\alpha _{n+2,p,q}`$ $`+(p+1)(p+2)\alpha _{n,p+2,q}+(q+1)(q+2)\alpha _{n,p,q+2}=0,`$
$`(n+1)(n+2)\alpha _{n+2,p,q}`$ $`(n+q)(p+2)\alpha _{n,p+2,q}+(q+1)(q+2)\alpha _{n,p,q+2}=0,`$
$`(n+1)(n+2)\alpha _{n+2,p,q}`$ $`+(p+1)(p+2)\alpha _{n,p+2,q}(n+p)(q+2)\alpha _{n,p,q+2}=0.`$
The determinant of the corresponding $`3\times 3`$-matrix is equal to $`2(n+2)(p+2)(q+2)(n+p+q+1)^2`$ and it follows that all the $`\alpha _{n,p,q}`$ are zero and, hence, $`\mathrm{deg}r<N`$, which gives a contradiction.
The argument for $`\stackrel{~}{S}_2`$ is entirely similar. ∎
Let $`(,)`$ be a non-degenerate symmetric bilinear form on a vector space $`V`$ of dimension greater than $`1`$. Define a ternary bracket on $`V`$ by
$$[a,b,c]=(a,c)b(b,c)a.$$
A straightforward verification shows that $`V`$ with this bracket satisfies all the axioms of a Lie triple system.
If $`I`$ is an ideal in $`V`$, $`[I,V,V]I`$, that is, $`(v,x)u(v,u)xI`$ for any $`xI`$ and $`u,vV`$. Hence, $`(v,x)uI`$ for any $`xI`$ and this means that $`I`$ is either trivial, or coincides with $`V`$. Therefore, $`V`$ is simple.
###### Proposition 26.
The L.t.s. $`V`$ satisfies Conjecture 3.
###### Proof.
Fix a basis $`\{x_k\}`$ for $`V`$, $`nk1`$, and let $`rU(V)`$ be homogeneous of degree greater than 1. The condition $`x_krrx_k=0`$ implies, by (13), that
$$\underset{i,j}{}[x_k,x_i,x_j]\frac{}{x_i}\frac{}{x_j}r=0,$$
that is,
$$\underset{i,j}{}((x_k,x_j)x_i(x_i,x_j)x_k)\frac{}{x_i}\frac{}{x_j}r=0.$$
Assuming that the basis $`\{x_k\}`$ is orthonormal, we get
$$x_k\underset{i}{}\frac{^2}{x_i^2}r=\underset{i}{}x_i\frac{}{x_i}\frac{}{x_k}r=(m1)\frac{}{x_k}r,$$
where $`m=\mathrm{deg}r`$. If $`\frac{}{x_k}r=0`$ for some $`k`$ it follows that $`_i\frac{^2}{x_i^2}r=0`$ and, hence, that $`\frac{}{x_k}r=0`$ for all $`k`$. In this case $`r`$ is a constant, so we can assume that $`\frac{}{x_k}r0`$ for all $`k`$ and that $`_i\frac{^2}{x_i^2}r0`$.
Let us write $`\psi `$ for $`_i\frac{^2}{x_i^2}r`$. We have
(14)
$$(m1)\frac{}{x_k}r=x_k\psi $$
and, hence,
$$(m1)x_k\frac{}{x_k}r=x_k^2\psi $$
and
$$m(m1)r=q\psi $$
with $`q=_{i=1}^nx_i^2`$. It follows that
$$m(m1)\frac{}{x_k}r=2x_k\psi +q\frac{}{x_k}\psi $$
which implies
$$(m2)x_k\psi =q\frac{}{x_k}\psi .$$
If $`m=2`$ this means that $`r`$ is a scalar multiple of $`q`$. If $`m2`$ we have that $`\psi =q\psi _0`$ with $`\psi _00`$, and $`m(m1)r=q^2\psi _0`$. It is readily seen that $`\psi _0`$ satisfies
$$(m4)x_k\psi _0=q\frac{}{x_k}\psi _0.$$
If $`m=4`$ this implies that $`r`$ is a scalar multiple of $`q^2`$; otherwise the above manipulations can be repeated. Eventually, this process has to stop and in the end we get that $`m=2l`$ and that, up to a multiplication by a scalar, $`r=q^l`$.
Now, (14) can be re-written as
$$(2l1)2x_klq^{l1}=x_k(2nlq^{l1}+4l(l1)q^{l1}).$$
This gives $`n=1`$ and it follows that $`x_krrx_k=0`$ cannot be satisfied for all $`k`$. ∎ |
warning/0506/cond-mat0506772.html | ar5iv | text | # Doped Mott Insulators are Insulators: Hole localization in the Cuprates
## Abstract
We demonstrate that a Mott insulator lightly doped with holes is still an insulator at low temperature even without disorder. Hole localization obtains because the chemical potential lies in a pseudogap which has a vanishing density of states at zero temperature. The energy scale for the pseudogap is set by the nearest-neighbour singlet-triplet splitting. As this energy scale vanishes if transitions, virtual or otherwise, to the upper Hubbard band are not permitted, the fundamental length scale in the pseudogap regime is the average distance between doubly occupied sites. Consequently, the pseudogap is tied to the non-commutativity of the two limits $`U\mathrm{}`$ ($`U`$ the on-site Coulomb repulsion) and $`L\mathrm{}`$ (the system size).
Hole doping a Mott insulator shiftssawatzky the chemical potential from the middle of the charge gap generated by the energy cost ($`U`$) for double occupancy to the top of the lower Hubbard band. Nominally, the density of states at the top of the lower Hubbard band is non-zero. Consequently, doped Mott insulators are expected to be conductors. However, doped Mott insulators such as the high temperature cuprate superconductors are well known to possess a pseudogapalloul ; norman ; timusk at the Fermi energy below some characteristic temperature, $`T^{}`$, that persists well into the superconducting dome. While a dip in the density of states is not sufficient to destroy the simple picture that a metallic state obtains upon light hole doping, certainly a vanishing density of the states at the Fermi level would be. The question arises: Does the density of states vanish at the chemical potential in the limit $`T0`$ in the underdoped cuprates or in lightly doped Mott insulators in general? The analysis presented here on the Hubbard model suggests the answer to this question is yes and hence lightly doped Mott insulators are, in fact, still insulators.
Experimental probes that shed light, either directly or indirectly, on the ultimate fate of the density of states at the chemical potential in doped Mott systems are of three types: transport, tunneling and angle-resolved photoemission (ARPES). Early transport measurements on La<sub>2-x</sub>Sr<sub>x</sub>CuO<sub>4+y</sub> revealedbirgeneau1 ; birgeneau2 that in the lightly-hole doped regime, the in-plane resistivity obeys the 3-d variable hopping form,
$`\rho (T)e^{(T_0/T)^\alpha }`$ (1)
with $`\alpha =1/4`$birgeneau1 or diverges logarithmically as $`\mathrm{ln}T_0/T`$birgeneau2 . In both cases, if localization is due to disorder (and hence extrinsic to Mott physics), an externally applied magnetic field should couple to the orbital motion and yield a negative magnetoresistance. While the magnetoresistance is negative, it is independent of the direction of the fieldbirgeneau2 , indicating that the localization mechanism is intrinsic and arises solely from spin scattering. In fact, extensive measurementsboeb1 ; boeb2 ; boeb3 ; boeb4 over the last 10 years indicate that once superconductivity is destroyed by the application of a large magnetic field, only two electrically distinct phases exist in the cuprates: 1) an insulator with a logarithmically diverging resistivity of the form $`\mathrm{ln}T_0/T`$ throughout the pseudogap region, $`x<x_\mathrm{c}`$, and 2) a Fermi liquid metal for $`x>x_\mathrm{c}`$, where $`x_cx_{\mathrm{opt}}`$, the optimal doping level. In the absence of a field, the most recentando in-plane transport data on untwinned crystals of YBa<sub>2</sub>Cu<sub>3</sub>O<sub>6.35</sub>, a composition right at the edge of the superconducting dome, corroborate the diverging $`\mathrm{ln}(T_0/T)`$ behaviour found in the high magnetic field limit for both $`\rho _a`$ and $`\rho _b`$. These authorsando conclude that the localization mechanism is independent of field and likely to be a consequence of an intrinsically insulating pseudogap at $`T=0`$. Scanning tunneling experimentsstm are consistent with the deepening of the pseudogap as the temperature is decreased. Finally, recent ARPES measurementsshen have detected a finite gap over the entire Brillouin zone, including along the $`d_{x^2y^2}`$ nodal line, in the normal state of La<sub>2-x</sub>Sr<sub>x</sub>CuO<sub>4</sub>, Ca<sub>2-x</sub>Na<sub>x</sub>CuO<sub>2</sub>Cl<sub>2</sub>, and Nd<sub>2-x</sub>Ce<sub>x</sub>CuO<sub>4</sub>. In Ca<sub>2-x</sub>Na<sub>x</sub>CuO<sub>2</sub>Cl<sub>2</sub>, the gap was observed to close at $`x=0.12`$. Similarly, in optimally doped Bi2212campu , the imaginary part of the self energy is momentum dependent but it remains non-zero even along the nodal directions. The ubiquity of a complete gap in the normal state of both electron and hole-doped cuprates prompted Shen, et. al.shen to conclude that gapped excitations in lightly doped Mott insulators is a generic feature.
Quite generally, a pseudogapefros ; altshuler is an example of an orthogonality catastrophemahan . Typically, orthogonality leads to vanishing of both the quasiparticle weight, $`Z`$, as well as the conductivity at $`T=0`$. In this case, we find that the orthogonality in a doped Mott insulator arises because hole transport is limited by the triplet-singlet energy gap. The length scale underlying this energy gap is the average separation between doubly occupied sites. As this length scale diverges in projected models but remains finite in the Hubbard model, the pseudogap is tied to a non-commutativity of $`U\mathrm{}`$ and $`L\mathrm{}`$. This lack of commutativity offers a possible explanation why all simulations thus far on the $`tJ`$ model find metallic transporthaule ; prelovsek near half-filling whereas for the Hubbard model, an insulating state obtains.
The starting point for our analysis is the Hubbard model. In this context, we have been refiningfmott a non-perturbative resolvent methodmm for calculating the single-particle spectral function $`A(𝐤,\omega )=\mathrm{Im}FT(\theta (tt^{})\{c_{i\sigma }(t),c_{j\sigma }^{}(t^{})\}`$ that is based on a self-consistent determination of the electron self-energy using the Hubbard operators. Here, $`c_{i\sigma }`$ is the electron annihilation operator and FT represents the frequency and momentum Fourier transform. In the spirit of cellular methodscell , the essence of our procedure is to expand the electron self-energy for the 2D lattice in terms of the resolvents for a small cluster. In our work, the eigenstates of a two-site cluster were used to expand the operators in the self-energy. As the self-energy can be written as a product of two operators, each of which can be centered on different lattice sites, a two-site expansion for each operator captures local correlations (albeit in a pair-wise fashion) over at most four lattice sites. Such a local expansion has been shown to yield a heat capacity of the 1D Hubbard systemfmott2 in excellent agreement with the Bethe ansatz as well as a pseudogapfmott in the 2D Hubbard model. Our emphasis here is on using the spectral function to calculate the conductivity. To obtain a direct link between the conductivity and the spectral function, we work with the non-crossing approximation
$`Re\sigma _{xx}(0+i\delta )`$ $`=`$ $`2\pi e^2{\displaystyle d^2k𝑑\omega ^{}(2t\mathrm{sin}k_x)^2}`$ (2)
$`\left({\displaystyle \frac{f(\omega ^{})}{\omega ^{}}}\right)\left[A(\omega ^{},k)\right]^2`$
to the Kubo formula for the conductivity where $`f(\omega )`$ is the Fermi distribution function. Although Eq. (2) is only approximate, as it does not include vertex corrections, we will show that our conclusions are independent of any approximation used to compute the conductivity. Shown in Fig. (1) is the resultant computation of the resistivity as a function of temperature for fillings of $`n=0.97`$, $`0.95`$, $`0.9`$, $`0.85`$, and $`n=0.8`$. At high temperatures the resistivity increases algebraically regardless of the filling. However, at low temperatues, a divergence in accord with Eq. (1) for fillings close to $`n=1`$ obtains. We will determine the crossover filling on general grounds later. In contrast, similar cluster treatments of the spectral functionhaule of the t-J model coupled with Eq. (2) find a metallic conductivity at all fillings, even arbitrarily close to half-filling. In fact, recentprelovsek exact diagonalization calculations on finite clusters confirm the inherent metallic behaviour at low temperatures, regardless of filling, in the t-J model. Metallic behaviour in the t-J model is consistent with the extensive numericalmischenko ; sorella ; dagotto and self-consistent Born calculationsborn which have found that a single hole is mobile in a quantum antiferromagnet described by the t-J model with a quasi-particle residue that scales as $`ZJ/t`$ where $`J=4t^2/U`$.
What then is the origin of the insulating state for the Hubbard model in the underdoped regime? The inset in Fig. (1) demonstrates that the density of states at the chemical potential plummets to zero exponentially as the temperature decreases. The conductivity, Eq. (2), is a product of the derivative of the Fermi distribution function and the spectral function. Because the former is peaked while the latter is zero at the chemical potential, the product necessarily vanishes leading to an insulating state. This cancellation persists to all orders of perturbation theory. Hence, the insulating state found here is not an artifact of the approximate form of Eq. (2); rather it arises simply because $`D(ϵ_F)=0`$ at $`T=0`$.
Because the electron self-energy is expanded in the level operatorsfmott for a two-site cluster, we can determine which local two-site correlations determine the physics of the vanishing of the density of states. The solid line in Fig. (2) illustrates clearly that the chemical potential lies in a local minimum in the single-particle density of states. This state of affairs obtains because nearest-neighbour singlet states (solid circles) and triplet (open squares) contribute to the density of states just below and above the chemical potential, respectively as shown in Fig. (2). Because the triplet and singlet are split by an energy $`J=4t^2/U`$, their contributions to the density of states cannot occur at the same energy. The density of states must have a dip which must constitute a real gap at $`T=0`$. The inset illustrates that precisely at the temperature (see Fig. (3)) at which the dip in the density of states obtains, the occupancy in the excited triplet states drops below that of the singlets. This definitively proves that it is the singlet-triplet excitation gap that limits hole transport in a doped Mott insulator.
Such a pseudogap can be thought of as a spin gapalloul as in the context of a spin liquidrvb . Also consistent with our finding here is the ferromagnetic polaron picturemanganites . However, neither experimental nor theoreticalaffleck work supports the ferromagnetic polaron model in the parameter range of the cuprates. In our simple picture that it is the local singlet-triplet splitting that gives rise to the pseudogap, we expect the corresponding gap arising from the orthogonality to be isotropic in momentum space. As illustrated by the inset in Fig. (3), the curvature of the density of states at the chemical potential is positive at each momentum indicating that all momenta contribute to the pseudogap, though with differing weights. This is consistent with the extensive ARPES study of Shen, et. al.shen . In the context of the cuprates, we propose that any anisotropyappears seen in the pseudogap is absent at $`T^{}`$ but arises at lower temperatures as a result of any ordering phenomenayazdani or pairing that might supervene on the pseudogap phase. In fact, othersjarrell have concluded recently based on cluster calculations on the Hubbard model that a pseudogap arises entirely from local correlations independently of any ordering or pair formation.
Two natural questions that arise from this work are 1) why do analogous cluster or exact diagonalization studies of the t-J model show no indication of localizationhaule ; prelovsek and 2) what sets the length scale for the energy gap. Both of these questions have the same answer. Without the triplet contribution, the pseudogap in Fig. (2) vanishes. However, the triplet contribution lies above the chemical potential and hence is part of the addition spectrum. The addition spectrumsawatzky of the low-energy spectral weight (LESW) is a sum of two distinct processes each involving spectral weight transfer between the upper and lower Hubbard bands: 1) a static part arising from state counting which grows as $`2x`$ but more importantly 2) a dynamical part that arises entirely from the hybridization. Since the triplet is present only when $`t0`$, the triplet contribution to the LESW is purely dynamical. In projected models in which double occupancy is eliminated at second order, the LESW scales exactly as $`2x`$sawatzky . Hence, the dynamical contribution to the spectral weight transfer is absent. However, the dynamical contribution to the addition part of the LESW can be treated perturbatively as first shown by Harris and Langeharris . Perturbation theory alone is insufficient to generate a gap in an excitation spectrum since the opening of a gap represents a phase transition. The essence of the problem is that as long as the insulating state is tied to the dynamical contribution to the spectral weight transfer between the upper and lower Hubbard bands, the length scale, $`\xi _{\mathrm{do}}`$, over which transport is governed by double occupancy must be finite. That is, the physics is sensitive to the order of limits of $`U\mathrm{}`$ and $`L\mathrm{}`$. Such non-commutativity signals a breakdown in perturbation theory as advocated previouslyphillips . $`U\mathrm{}`$, $`L\mathrm{}`$ results in $`\xi _{\mathrm{do}}>L`$, metallic transport. In the reverse order of limits, $`\xi _{\mathrm{do}}<L`$ and localization obtains provided that the $`n_h\xi _{\mathrm{do}}^2<L^2`$, $`n_h=x(L/a)^2`$ the number of holes. $`n_h\xi _{\mathrm{do}}^2=L^2`$ defines the percolation limit. By calculating the percentage of doubly occupied sites, we obtained $`\xi _{\mathrm{do}}`$ numerically and plotted the $`T^{}`$-line, $`J(1cx(\xi _{\mathrm{do}}/a)^2)`$, in Fig. (3). The agreement of this phenomenological fit with the resisitivity data in which a metallic state obtains at $`x=0.1`$ and the crossing in Fig. (2) lend credence to our assertion that $`\xi _{\mathrm{do}}`$ is the relevant length scale for the pseudogap. Finally, the scaling form
$`ZL^{(t/U)^p}p>0`$ (3)
for the one-hole quasiparticle weight lays plain that the discrepancy between the $`tJ`$mischenko ; dagotto ; born and Hubbardsorella results is one of lack of commutativity. In the $`tJ`$ model (no double occupancy), $`U\mathrm{}`$, $`L\mathrm{}`$ and $`Z`$ remains finite. In the reverse order of limits (Hubbard model), $`Z`$ vanishes.
Indeed, other proposals for hole localization exist. Some have argued that in the t-J model, a hole creates a phase stringweng . However, such an exotic state is not borne out by extensive numerical simulations on the t-J modelmischenko . In the spin-fermion model, selective gapping occurs at hot spots indicated by the intersection of the Fermi surface arcs with the reduced diamond-shaped AF Brillouin zoneschmalian whereas in the spin-bag modelbag , a gap occurs only along the $`(\pi ,\pi )`$ direction. Neither of these models, however, possesses the strong correlations intrinsic to the doped Mott state.
To conclude, our proposal that an orthogonality between the singlet and triplet states necessarily requires a finite length scale over which transport is governed by the distance between double occupancies implies that $`U\mathrm{}`$ and $`L\mathrm{}`$ do not commute. The emergence of such a finite length scale in the transport properties offers a possible resolution of the breakdown of the one-parameter scaling pictureone for quantum criticality in the cuprates. Finally, experimentsdiamag demonstrating that diamagnetism in the pseudogap phase does not persist all the way to $`T^{}`$ proves that pair fluctuations alone cannot account for the pseudogap. As advocated here, the pseudogap arises from Mottness and any relationship between orderingyazdani or pairing and the pseudogap is one of supervenience.
###### Acknowledgements.
We thank the NSF, Grant No. DMR-0305864, Duncan Haldane for a key discussion on the $`U\mathrm{}`$ and $`L\mathrm{}`$ limits and also T. Stanescu and A. Yazdani. |
warning/0506/astro-ph0506728.html | ar5iv | text | # Photometric Periodicities of Be/X-ray Pulsars in the Small Magellanic Cloud1footnote 11footnote 1This paper utilizes public domain data obtained by the MACHO Project, jointly funded by the US Department of Energy through the University of California, Lawrence Livermore National Laboratory under contract No. W-7405-Eng-48, by the National Science Foundation through the Center for Particle Astrophysics of the University of California under coopertative agreement AST-8809616, and by the Mount Stromlo and Siding Spring Observatory, part of the Australian National University.
## 1 Introduction
Very luminous, high-mass X-ray binaries (HMXB) are generally divided into two groups: 1) those containing a Roche-lobe filling supergiant which transfers mass onto a neutron star or black-hole companion, and 2) wider binaries containing a Be star and a neutron star which accretes material from the equatorial disk of the primary. In the latter systems the compact star is often found to be an X-ray pulsar. In the Galaxy the HMXB are about equally divided between these two broad classes, but in the Magellanic Clouds the Be/pulsar systems make up the majority of HMXB. In the SMC we recognize only one supergiant system (SMC X-1), while several dozen X-ray pulsars with Be-star companions are known (e.g. Haberl & Sasaki 2000; Coe et al. 2002).
Only a few of the Be/neutron star systems have been studied in detail, but from these it appears the stars are in fairly wide but eccentric orbits. Hence the two stars interact primarily near periastron passage when the neutron star may enter the extended equatorial disk of the Be star and experience increased accretion of gas. This gives rise to both X-ray and optical outbursts from which the orbital period can be determined. Corbet (1984) and others (e.g. in’t Zand et al. 2001; Charles & Coe 2004) have shown that the pulse and orbital periods in such Be/X-ray systems are related. This correlation was explained by Waters & van Kerkwijk (1989) as resulting from the slow equatorial Be winds which affect the equilibrium spin period of the pulsar.
A notable example of periodic orbital outbursts is seen in the LMC source A0535$``$668, where the eccentricity is $``$0.7. This system shows extreme changes in spectrum and brightness at all wavelengths near periastron (e.g. Charles et al. 1983; Hutchings et al. 1985). For many of the other Be/X-ray systems we see less dramatic but clear periodic outbursts (e.g. AX J0049.4$``$7323, Cowley & Schmidtke 2003), although the orbital eccentricity has not yet been determined for almost all of these systems. Schmidtke et al. (2004) identified new orbital periods in 7 Magellanic Cloud Be/X-ray systems based on longterm MACHO and OGLE-II photometric data. In addition to orbital outbursts, they found quasi-periodic variations (QPV) which they suggested were due to reverberations in the Be star’s equatorial disk following penetration by the neutron star. Generally the QPV have timescales of days. The QPV periods generally differ by a few percent following each orbital interaction.
As pointed out by Schmidtke et al. (2004, and references therein) Be/neutron star binaries often show longterm light curves with considerable variability. Some exhibit prominent ‘swoops’ or irregular outbursts (several examples are shown in Fig. 4 of Coe et al. 2005), while others show a lot of scatter in the light curve. For a few, their periodic behavior is clearly visible without any analysis (e.g. RX J0049.4$``$7323, Cowley & Schmidtke 2003). However, for most systems a full analysis of the photometry is needed to identify any periodic signal.
In the present paper we have investigated 7 more SMC systems in an effort to better understand the variability seen in Be/X-ray-pulsar systems. The particular systems were selected on the basis of several criteria. First, they have well determined positions, and hence we are confident in the optical identifications. Second, the sources had to be included in either the MACHO or OGLE surveys. Third, we chose sources where the variability exceeds the photometric errors. Additionally, the OGLE sources are all in their “Difference Image Analysis” (DIA) catalogue”, so they had already been selected as variables.
In addition to orbital outbursts and QPV, we have discovered that some of these Be stars also show nonradial pulsations (NRP) with timescales of hours to over a day. In retrospect, this is not surprising, since it is known that some single Be stars show such variability (e.g. Balona, Sterken, & Manfroid 1991; Balona 1992; Balona & James 2002; Percy, Harlow, & Wu 2004).
## 2 Analysis of Optical Data from the MACHO and OGLE-II Projects
Our photometric data come from two large surveys which allow public access to their data: the OGLE-II survey (Udalski, Kubiak, & Szymanski 1997; Zebrun et al. 2001) and the MACHO survey. OGLE-II provides photometry in $`I`$ for variables in their DIA catalogue. The MACHO site gives ‘blue’ and ‘red’ instrumental magnitudes for all stars in their fields, and these data can be transformed to standard $`V`$ and $`R`$ colors (Alcock et al. 1999). (Note: For users looking at the web display of ‘calibrated’ MACHO light curves for SMC sources, one must add 0.75 mag to the scale shown, since only half of the actual integration time (300 s versus 600 s) is used when generating these quick-look light curves.)
Table 1 gives the equinox 2000 MACHO position of the optical stars, the mean $`R`$ magnitudes, the mean $`VR`$ colors, the MACHO and OGLE-II identification numbers, and the catalogue numbers from “H$`\alpha `$-bright Stars in the SMC” (Meyssonnier & Azzopardi 1993; hereafter called MA93) for the systems we studied.
Each data set ($`V`$, $`R`$, $`I`$) was flattened, or pre-whitened, to remove longterm trends. This was done by fitting a low-order polynomial to long stretches of data in each color, excluding extreme outlying points or sections of the light curve where unrepeated sudden changes occurred. The resulting “detrended” magnitudes ($`V^{}`$, $`R^{}`$, $`I^{}`$) were then analyzed. In some cases we separated the MACHO data by east-pier and west-pier designations prior to detrending in order to remove any systematic differences between the two instrumental configurations.
Using the method described by Horne & Baliunas (1986), we searched for periodic behavior in the detrended light curves. The resulting periodograms are shown for each system discussed below. In cases where the periodic light curves are very non-sinusoidal, this method of analysis may display only low power at the true period, since it searches for sinusoidal variations. Usually aliases of the true period are stronger, as is the case for several of the systems studied in this paper. In these cases the period is better revealed using the phase dispersion minimization (PDM) technique described by Stellingwerf (1978). The PDM method can identify periodicities in light curves having an arbitrary but repeatable shape, such as the recurring outbursts seen in some of the systems described here. More details are given with each source where the PDM method has been used. Typically, we searched for periods in the range 0.25-1000 days, covering a minimum of 10<sup>4</sup> test frequencies.
For a given periodogram, the significance of its peaks was explicitly assessed by creating 1000 data sets, with randomized magnitudes and dates. From a simple tally of the strongest power peak within each of these sets, we determined the 99% confidence level. That is, in only 1% of the random data sets did the peak power exceed the stated 99% level. Nearly all of the periods found in this study are significant at or above this value. This assessment assumes that the signal is strictly sinusoidal and coherent. If these conditions do not hold, such as in QPV or NRP, then the observed power can be somewhat less than the ideal case, so that the significance of the peak appears to be lower.
Most of the figures display folded light curves, using the “detrended” magnitudes. Hence, these light curves no longer contain information about the brightness of the system. However, the true magnitudes are plotted in the longterm light curves.
## 3 Individual Be/X-ray Pulsar Systems
### 3.1 XMMU J004723.7$``$731226 = RX J0047.3$``$7313
This 263.6 s X-ray pulsar was discovered by Ueno et al. (2004) but was also studied by Haberl & Pietsch (2004). Because of its very good XMM-Newton position, we can confidently identify it with MACHO 212.15792.77. It also appears to be MA93#172, confirming the primary as an emission-line B star.
The longterm MACHO $`V`$ light curve is shown in the upper panel of Fig. 1. Both the MACHO and OGLE-II data show a general scatter of $``$0.1 mag superimposed on a slow brightening trend. This system is one of the cases where our standard analysis reveals aliases rather than the fundamental period. The periodogram from the $`I^{}`$ data shows power at 12.3 d and 9.8 d, but light curves folded on these periods are very scattered. The PDM analysis of the $`I^{}`$ data shows that the fundamental period is P=49.1$`\pm `$0.2 days (third panel of Fig. 1), which we assume is the orbital period. However, both the periodogram and the PDM variance for the $`V^{}`$ and $`R^{}`$ data show no significant periods. In similar Be/X-ray systems the orbital outbursts are usually strongest in $`I`$, and hence the lack of periodic variation in $`V`$ or $`R`$ is not surprising. Edge et al. (2005a) report P=48.8$`\pm `$0.6 d, also from OGLE $`I`$ data, although details of their analysis are not yet available.
The folded orbital $`I^{}`$ and $`V^{}`$ light curves are shown in the bottom panel of Fig. 1. In $`I^{}`$ there is a small outburst lasting for $``$0.1P, with an amplitude of $``$0.02 mag. Examining different segments of the data reveals that from cycle to cycle the shape of the light curve is similar, but some outbursts are stronger than others. For a short time before and after each outburst the source appears to be $``$0.01 mag fainter than the mean. There is no corresponding outburst in $`V^{}`$ or $`R^{}`$.
### 3.2 RX J0049.1$``$7250 = AX J0049$``$729
X-ray pulses at 74.7 s from RX J0049.1$``$7250 were discovered by Corbet et al. (1998) and confirmed by Yokogawa & Koyama (1998a). The optical counterpart is identified with the emission line Star #1 of Stevens, Coe, & Buckley (1999). The system has a mean magnitude of $`R`$16.9, and its longterm light curve shows variability of $``$0.2 mag (see Fig. 2).
MACHO light curves in $`V`$ and $`R`$ were analyzed in two segments (A and B, as marked on Fig. 2). Both the OGLE-II and MACHO photometric data show a highly significant periodicity at 33.4$`\pm 0.4`$ days, with its detection being well over the 99% confidence level. This period is likely to be the orbital period, and it is in good agreement with the relation between orbital and pulse periods first recognized by Corbet (1984). The second panel of Fig. 2 shows the periodogram of $`R^{}`$ data taken from segments A and B. The $`R^{}`$ light curve folded on P=33.4 days, is shown in the third panel. In $`V^{}`$ light, the 33-day period is present but weaker.
The periodogram also reveals several significant peaks with periods in the range of 2.39-2.40 days which we identify as quasi-periodic variations of the Be star’s disk. The QPV have the dominant power in $`V^{}`$ where the 33-day signal is weaker. The $`R^{}`$ and $`V^{}`$ data from Seg A, folded on the QPV 2.4-day period, are plotted in the bottom panel of Fig. 2.
The longterm $`I`$ light curve for RX J0049.1$``$7250 is shown in the upper panel of Fig. 3. Note the super-outburst near MJD 50760. The middle panel displays the periodogram from $`I^{}`$ data, with the observations near the super-outburst excluded. The 33.4-day period is prominent. The bottom panel shows the $`I`$ data folded on this period, with the super-outburst superimposed in a different symbol (open circles). The mean amplitude of the 33-day variation is $`\mathrm{\Delta }`$$`I`$0.03 mag. The super-outburst is in phase with the other periodic brightenings – it just happened to be much larger, perhaps due to denser material surrounding the Be star at the time of the interaction.
We note that Coe & Orosz (2000) reported that they found no periods in the OGLE-II data in the range 1-50 days. Laycock et al. (2004) suggested a possible orbital period of 642$`\pm `$59 days based on the intervals between three strong X-ray outbursts observed with RXTE. However, with more X-ray data Galache (2005, private communication) has found P<sub>X</sub>$``$65.1 days, which may be double the 33.4-day optical period identified here. This suggests that a strong X-ray outburst is not seen at every periastron passage.
### 3.3 XTE J0052$``$725
The position of this 82.4 s pulsar was determined using $`Chandra`$ data by Edge et al. (2003). The longterm $`R`$ light curve shows the mean magnitude was fairly constant until MJD 50800, with a scatter of $``$0.1 mag. After that date the source slowly declined by $``$0.3 mag in $`B`$, $`V`$, and $`I`$. Both $`R`$ and $`I`$ longterm light curves are plotted in Fig. 4. Coe et al. (2005; their Fig. 4) using both OGLE-II and -III data show that the source brightened slightly between MJD 51500 and 52000, but then fell rapidly by at least 0.6 mag.
The MACHO $`V^{}`$ and $`R^{}`$ light curves were analyzed in four time segments (A-D), as shown in Fig. 4. A strong periodic variation, with P=1.328 days, is present in both MACHO colors. The $`R^{}`$ power spectrum for segments A through C is shown in the middle panel of Fig. 4; the primary peak and its aliases are marked. Examining the individual segments, the power steadily dropped until there was virtually no periodic signal in Seg D (see bottom panel of Fig. 4). The period remained essentially constant while the variation was present.
As expected, the $`I^{}`$ data which overlaps Seg D (from MJD $``$50500 to $``$51500) also shows no periodic signal. However, in spite of the small amount of $`I`$ data in Seg E, a power spectrum analysis reveals the original period returned during this segment, although its phasing differed by about a quarter of a cycle. Such a short period with variable amplitude suggests that it is due to nonradial pulsations of the Be star. This is the first time NRP have been identified in a Be/X-ray binary, although similar short periods are observed in some single Be stars (e.g. Percy et al. 2004; Balona & James 2002; Stefl & Balona 1996).
$`R^{}`$ and $`V^{}`$ light curves from Seg A, B, and C, folded on P=1.328 days, are shown in Fig. 5. In a given segment, both colors have similar amplitudes. The greatest amplitude occurred in Seg A ($`\mathrm{\Delta }`$m$``$0.04 mag). In subsequent segments the amplitude decreased ($`\mathrm{\Delta }`$m$``$0.02-0.03 mag in Seg B; $`\mathrm{\Delta }`$m$``$0.01 mag in Seg C) until there was no obvious periodic variation in Seg D. The $`I^{}`$ light curve from Seg E, folded on the mean period from Seg A and B, is shown in the bottom panel of Fig. 5. The amplitude is similar to that in Seg A, but the phasing differs as noted above.
We searched for but were unable to find any photometric signature of an orbital period up to P=1000 days.
### 3.4 CXOU J005455.6$``$724510 = RX J0054.9$``$7245 = AX J0054.8$``$7244
CXOU J005455.6$``$724510 appears to have been independently discovered as a $``$500 s pulsar by Edge et al. (2004a) using $`Chandra`$ archival data and by Haberl et al. (2004a) from XMM-Newton data. Edge et al. (2004b) give the pulse period as 503.5$`\pm `$6.7 s, while Haberl et al. (2004b) find P$`{}_{pulse}{}^{}=499.2\pm `$0.7 s.
The highly accurate $`Chandra`$ position allows one to confidently identify it as the emission-line star MA93#809 and MACHO 207.16254.16. Although the field is covered by OGLE-II data, the star is not listed as a variable in the DIA catalogue. This is probably due to the fact that the amplitude of the photometric variation is quite small.
The longterm $`R`$ light curve of CXOU J005455.6$``$724510 is plotted in Fig. 6. Both a periodogram and the variance from PDM analysis were used to study the $`R^{}`$ and $`V^{}`$ data. The $`R^{}`$ data show significant power at P=273$`\pm `$6 days (see middle panel in Fig. 6; the 273-day peak is well over the 99% confidence level). We infer that this is the orbital period. The peak near 0.002 d<sup>-1</sup> is not seen in the PDM analysis and therefore is likely to be spurious. The $`R^{}`$ light curve folded on P=273 days shows a small outburst ($`\mathrm{\Delta }`$$`R`$$``$0.02 mag) with a rapid rise and slower decline (see bottom panel). The $`V^{}`$ data show no significant periods, but as pointed out above for XMMU J004723.7$``$731226, orbital outbursts tend to be strongest in $`I`$, weaker in $`R`$, and may not be present in $`V`$. Folding the $`V^{}`$ data on 273 days shows that if there is any outburst in this color, its amplitude is $`<`$0.01 mag. A recent announcement by Edge et al. (2005b) shows they have independently found a similar period (P=268$`\pm `$1.4 d) using OGLE and MACHO data.
The optical data were also examined for shorter periods, looking for either QPV or NRP. There do not appear to be any significant periods less than a day, but there is a weak signal at $``$5.3 days, suggesting that QPV may be present after some outbursts. This period is present in both the full data set and when we isolate data taken between the orbital outbursts. Its significance is at the 98% confidence level. However, as discussed in Section 2, the power of QPV signals may be weaker in the full data set since there is not a single, stable period. In Fig. 7 we plot the $`R^{}`$ light curve folded on the 5.3-day period.
### 3.5 RX J0054.9$``$7226 = XTE J0055$``$724
The 58.9s pulsar XTE J0055$``$724 was discovered by Marshall et al. (1998). It is identified with MACHO 207.16259.23 and the emission line star MA93#810. Its longterm $`R`$ light curve is shown in the upper panel of Fig. 8. From a subset of OGLE-II data, Coe & Orosz (2002) found a peak in the frequency spectrum corresponding to a period of 14.26 days, but they say no obvious light curve was present when they folded the data on this period.
We have analyzed all of the OGLE-II data as well as the $`V`$ and $`R`$ photometry from the MACHO project. Periodograms for detrended $`R`$ and $`I`$ data show power near 15, 20, and 30 days, but these turn out to be aliases of the true period. Because of the presence of multiple periods, we also examined portions of the data in more detail – that is, divided into smaller time segments and by MACHO pier orientation. Each subsection of data gave somewhat different power for periods near 15, 20, and 30 days. However, analysis using the PDM method shows the fundamental (orbital) period is 60.2$`\pm `$0.8 days, as is seen in both $`R^{}`$ and $`I^{}`$ data in Fig. 8. The non-sinusoidal shape of this light curve produces no prominent power in the periodogram at the fundamental period. This is similar to the behavior seen in XMMU J004723.7$``$731226 (discussed above) where the aliases are simple fractions of the orbital period.
The folded $`R^{}`$ and $`I^{}`$ light curves can been seen in Fig. 9. They have a peculiar shape, which looks like a sinusoidal variation with a superimposed outburst near phase zero. The outburst itself barely rises above the mean light level. The light curve resembles that found for XMMU J004723.7$``$731226 (see Fig. 1), although the ‘outburst’ here is less pronounced.
Using RXTE data, Laycock et al. (2004) proposed a period of P<sub>X</sub>=123$`\pm `$1 days based on the spacing between four X-ray outbursts. We note that this is about double the optically determined period. We also searched for short periods that might arise from either QPV or from NRP, but we found nothing significant.
### 3.6 XMMU J005517.9$``$723853 = RX J0055.2$``$7238
RX J0055.2$``$7238 is a 701.6 s X-ray pulsar (Haberl et al. 2000; Haberl et al. 2004) which is identified with MACHO 207.16313.35. Its longterm light curve is shown in Fig. 10. Clearly there is shortterm variability around a mean magnitude of $`V`$16.1. The slight downward trend in both $`V`$ and $`R`$ was removed before further analysis of the data.
The detrended photometric data reveal a pronounced period at P=0.28 days, much too short to be either an orbital period or even the Be star’s rotational period. The variation is approximately sinusoidal with amplitudes of $`\mathrm{\Delta }`$$`V`$0.03 and $`\mathrm{\Delta }`$$`R`$0.02 mag. This short period is most likely due to nonradial pulsations of the type seen in some single Be stars (e.g. Balona & James 2002). This is the second example (see XTE J0052$``$725 above) of NRP discovered in this study.
The periodogram from the $`V^{}`$ data is plotted in Fig. 10, where aliases of the 0.28-day period are also seen. The same period is found using $`R^{}`$ data. The bottom panel shows the $`V^{}`$ and $`R^{}`$ light curves folded on this very short period. Porter & Rivinius (2003; also see other references therein) in their review of classical Be stars report that in single Be stars, nonradial pulsations are more common in early spectral types. Haberl et al. (2004) found that the color of RX J0055.2$``$7238 implies the primary star has a very early spectral type (O9V), making it consistent with the behavior of single Be stars.
When the MACHO data are subdivided into 4 time segments (A-D), we find that the period changed with time. Expanded periodograms for each segment are shown in Fig. 11. In the lower panel we plot period versus time, showing that the period decreased from 0.28484 days to 0.28470 days at a mean rate of $`2.6`$ s yr<sup>-1</sup>. NRP in Be stars are well known to exhibit changes in periods, amplitudes, and light curve shapes similar to those observed here (e.g. Balona, Sterken, & Manfroid 1991).
Analysis of the $`R^{}`$ data taken before and after removal of the 0.28-day variation shows weak power at P$``$412$`\pm `$4 days in both the periodogram and the PDM. The power spectrum in Fig. 12 shows the 412-day peak has a confidence level below 90%. If real, we assume this is the orbital period, which would be in reasonable agreement with the Corbet (1984) P(pulse)/P(orbit) relation. The 412-day $`R^{}`$ light curve shows $``$0.01 mag variation (bottom panel of Fig. 12), similar to the behavior found in some other Be/X-ray systems. If the orbit is not very eccentric it is possible that there is little interaction between the neutron star and the Be star’s equatorial disk, resulting in only a very small change in system brightness.
### 3.7 Reanalysis of RX J0050.7$``$7316 = AX J0051$``$733 = DZ Tuc
The discovery of this 323 s pulsar was announced by Yokogawa & Koyama (1998b) and discussed in more detail by Imanishi et al. (1999). The optical counterpart was identified as a Be star by Cowley et al. (1997) in their study of X-ray sources in the Magellanic Clouds, and it is coincident with MA93#387. The colors ($`BV=0.03`$ and $`UB=0.95`$) indicate that the primary is a very early B or late O star.
Using OGLE-II data, Coe & Orosz (2000) found a period of 0.708 days, as had been suggested by Cook (1998), which they interpreted as half the orbital period in a close binary. Further analysis using a subset of MACHO data as well as OGLE-II data was carried out by Coe et al. (2002) showing that both the period and amplitude change with time. If the period change is linear, it amounts to $``$13.5 s yr<sup>-1</sup>, with the $`R`$ amplitude changing by $``$40%. In spite of the extreme nature of these parameters, they concluded that the system was likely to contain a very close binary with P<sub>orb</sub>$``$1.4 days.
We have reanalyzed the photometry for this system, now using the complete MACHO data set which adds more than 3 years of data. The longterm $`V`$ light curve is shown in Fig. 13. The amplitude clearly changes with time, with the least variability occurring around MJD 50000. We also found that there are slightly different mean magnitude levels in the MACHO data depending on whether the telescope was east or west of the pier. Thus, we flattened the data from these two configurations separately and then combined them to search for periodicities.
Our period analysis was done first using the entire data set and then subdividing the data into time segments of $``$200 days, some of which are indicated on Fig. 13. In the middle panel of Fig. 13 we show the power spectrum for several sample segments, demonstrating that the period changes. Light curves, folded on the $``$0.7-day period, are shown for two of these segments to emphasize how much the amplitude varies. Within a given time segment, there is little or no color variation through the 0.7-day period.
In Fig. 14 we plot the photometric period and amplitude versus date, using only segments which contain more than 30 observations. The character of the variation, with both period and amplitude changing with time, suggests that RX J0050.7$``$7316 is another nonradial pulsator, as was found above for XTE J0052$``$725 and XMMU J005517.9$``$723853. Some single Be stars are known to show pulsation periods near 0.7 days (e.g. Balona et al. 1991, Balona 1992). The 1.4-day binary interpretation always required special circumstances, such as rapid mass transfer, to explain all of the observations. Instead, interpreting the variations as due to NRP seems more probable to us, since two other systems with very short periods have been found. In addition, the early spectral type of this Be star suggests it might be likely to show NRP, as discussed above.
In order to search for a longer binary period we removed the 0.7-day variations and examined the residuals. Although we searched up to P=1000 days, there was no power strong enough at any frequency to be convincing evidence of orbital outbursts.
Imanishi et al. (1999) suggested their X-ray data were consistent with a 185-day period, but they stressed that further monitoring would be needed. Laycock et al. (2004) proposed a possible 108$`\pm `$18 day X-ray period. We find nothing in the optical data to support either of these suggested periods. Perhaps the orbit is fairly circular so that the neutron star does not come close enough to the Be star to cause an optical outburst, hence leaving the primary to behave like an isolated star.
## 4 Discussion and Summary
In this paper we have investigated the longterm light curves of 7 Be/X-ray pulsars in the Small Magellanic Cloud. We have shown that multiple types of optical variability are present including longterm irregular changes, orbital outbursts, quasi-periodic variations of the disk, and nonradial pulsations of the Be star.
For 5 of the systems studied, we have found low amplitude periodic outbursts which we interpret as being due to the interaction of the neutron star with the Be star’s disk once during each orbit, probably near periastron passage. Corbet (1984) and others have shown that the X-ray pulsation period and orbital period are correlated in Be/X-ray binaries. In Fig. 15 we plot a “Corbet diagram”, including the systems in this study as well as other Be/X-ray binaries (both in the Galaxy and the Magellanic Clouds) with published orbital periods.
For two systems (RX J0050.7$``$7316 and XTE J0052$``$725) we did not find any signature of the orbital period in the light curve. However, for both of these systems we discovered that NRP are present. These nonradial pulsators may either be in sufficiently wide orbits that the Be star behaves like a single star or the system could have a low eccentricity so there is little interaction between the stars. Although such pulsations are known to be present in some single Be stars, this is the first time they have been identified in Be/X-ray binaries. In addition, we found that XMMU J005517.7$``$723853 shows both strong NRP and a small outburst every $``$412 days (probably the orbital period).
RX J0049.1$``$7250 and CXOU J005455.6$``$724510 show both orbital outbursts and quasi-periodic variations of the type previously recognized by Schmidtke et al. (2004). The strong signature of the 33-day orbital period in RX J0049.1$``$7250 suggests that its orbit may be significantly eccentric.
In summary, we have discovered that nonradial pulsations are present in some Be/X-ray systems. We have also determined orbital periods for 5 systems and found 2 which show quasi-periodic variations. The Be/X-ray binaries may be divided into at least two groups. The systems in which the Be star has NRP either show no orbital outbursts or very weak ones, suggesting that they are in wide or nearly circular orbits (or both) with little interaction between the two stars. The systems in which the orbital interaction is strong (outbursts present) often exhibit QPV following the outbursts. None of the studied binaries with QPV also shows nonradial pulsations. The remaining two systems in which the orbital signature is weak show neither QPV nor NRP.
To further investigate these types of behavior, we reexamined the photometry of the 4 Be/X-ray binaries studied previously by Schmidtke et al. (2004) which display both orbital and QPV variations. We searched for very short periods which would indicated the presence of nonradial pulsation, but none was found. Hence these binaries behave like the ones in the present study, with QPV and NRP not occurring in the same system.
We thank Robin Corbet, Malcomb Coe, and Jose Galache for helpful information about some of these systems. We also thank the anonymous referee for useful suggestions which improved the paper. |
warning/0506/physics0506112.html | ar5iv | text | # Geometrical aspects of first-order optical systems
## 1 Introduction
Matrix methods offer the great advantage of simplifying the presentation of linear models and clarifying the common features and interconnections of distinct branches of physics . Modern optics is not an exception and a wealth of input-output relations can be compactly expressed by a single matrix . For example, the well-known $`2\times 2`$ ray-transfer matrix, which belongs to the realm of paraxial ray optics, predicts with almost perfect accuracy the behavior of a Gaussian beam.
In this respect, we note that there is a wide family of beams (including Gaussian Schell-model fields, which have received particular attention ) for which a complex parameter can be defined such that, under the action of first-order systems, it is transformed according to the famous Kogelnik $`ABCD`$ law . This is the reason why they are so easy to handle. This simplicity, together with the practical importance that these beams have for laser systems, explain the abundant literature on this topic .
The algebraic basis for understanding the transformation properties of such beams is twofold: the ray-transfer matrix of any first-order system is an element of the group SL(2, $``$ and the complex beam parameter changes according to a bilinear (or Möbius) transformation .
The nature of these results seems to call for a geometrical interpretation. The interaction between physics and geometry has a long and fruitful story, a unique example is Einstein theory of relativity. The goal of this paper is precisely to provide such a geometrical basis, which should be relevant to properly approach this subject.
The material of this paper is organized as follows. In section 2 we include a brief review of the transformation properties of Gaussian beams by first-order systems, introducing a complex parameter $`Q`$ to describe the different states as points of the hyperbolic plane. The action of the system in terms of $`Q`$ is then given by a bilinear transformation, which is characterized through the points that it leaves invariant. From this viewpoint the three basic isometries of this hyperbolic plane (i.e., transformations that preserve the distance), namely, rotations, translations, and parallel displacements, appear linked to the fact that the trace of the ray-transfer matrix has a magnitude lesser than, greater than, or equal to 2, respectively.
In section 3 we present a mapping that transforms the hyperbolic plane into the unit disc (which is the Poincaré model of the hyperbolic geometry) and we proceed to study the corresponding motions in this disc. Finally, as a direct application, in section 4 we treat the case of periodic systems, which are the basis for optical resonators, providing an alternative explanation of the standard stability condition.
We emphasize that this geometrical scenario does not offer any advantage in terms of computational efficiency. Apart from its undeniable beauty, its benefit lies in gaining insights into the qualitative behaviour of the beam evolution.
## 2 First-order systems as transformations in the hyperbolic plane $``$
We consider the paraxial propagation of light through axially symmetric systems, containing no tilted or misaligned elements. The reader interested in further details should consult the extensive work of Simon and Mukunda . We take a Cartesian coordinate system whose $`Z`$ axis is along the axis of the optical system and represent a ray at a plane $`z`$ by the transverse position vector $`x(z)`$ (which can be chosen in the meridional plane) and by the momentum $`p(z)=n(z)dx/dz`$ . Here $`n(z)`$ is the refractive index and $`dx/dz`$ is the direction of the ray through $`z`$.
At the level of ray optics, a first-order system changes the ray parameters by the simple transformation
$$\left(\begin{array}{c}x^{}\\ p^{}\end{array}\right)=𝐌\left(\begin{array}{c}x\\ p\end{array}\right),$$
(2.1)
where the primed and unprimed variables refer to the output and input planes, respectively, and $`𝐌`$ is the ray-transfer matrix that must satisfy the condition
$$𝐌=\left(\begin{array}{cc}A& B\\ C& D\end{array}\right),det𝐌=ADBC=1,$$
(2.2)
which means that $`𝐌`$ is an element of the group SL(2, $``$) of real unimodular $`2\times 2`$ matrices.
When one goes to paraxial-wave optics, the beams are described in the Hilbert space $`L^2`$ of complex-valued square-integrable wave-amplitude functions $`\psi (x)`$. The classical phase-space variables $`x`$ and $`p`$ are now promoted to self-adjoint operators by the procedure of wavization , which is quite similar to the quantization of position and momentum in quantum mechanics.
We are interested in the action of a ray-transfer matrix on time-stationary fields. We can then focus the analysis on a fixed frequency $`\omega `$, which we shall omit henceforth. Moreover, to deal with partially coherent beams we specify the field not by its amplitude, but by its cross-spectral density. The latter is defined in terms of the former as
$$\mathrm{\Gamma }(x_1,x_2)=\psi ^{}(x_1)\psi (x_2),$$
(2.3)
where the angular brackets denote ensemble averages.
There is a wide family of beams, known as Schell-model fields, for which the cross-spectral density (2.3) factors in the form
$$\mathrm{\Gamma }(x_1,x_2)=[I(x_1)I(x_2)]^{1/2}\mu (x_1x_2),$$
(2.4)
where $`I`$ is the intensity distribution and $`\mu `$ is the normalized degree of coherence, which is translationally invariant. When these two fundamental quantities are Gaussians
$`I(x)`$ $`=`$ $`{\displaystyle \frac{}{\sqrt{2\pi }\sigma _I}}\mathrm{exp}\left({\displaystyle \frac{x^2}{2\sigma _I^2}}\right),`$
$`\mu (x)`$ $`=`$ $`\mathrm{exp}\left({\displaystyle \frac{x^2}{2\sigma _\mu ^2}}\right),`$
the beam is said to be a Gaussian Schell model (GSM). Here $``$ is a constant independent of $`x`$ that can be identified with the total irradiance. Clearly, $`\sigma _I`$ and $`\sigma _\mu `$ are, respectively, the effective beam width and the transverse coherence length. Other well-known families of Gaussian fields are special cases of these GSM fields. When $`\sigma _\mu \sigma _I`$ we have the Gaussian quasihomogeneous field, and the coherent Gaussian field is obtained when $`\sigma _\mu \mathrm{}`$. In any case, the crucial point for our purposes is the observation that for GSM fields one can define a complex parameter $`Q`$
$$Q=\frac{1}{R}+i\frac{1}{k\sigma _I\delta },$$
(2.6)
where
$$\frac{1}{\delta ^2}=\frac{1}{\sigma _\mu ^2}+\frac{1}{(2\sigma _I)^2},$$
(2.7)
and $`R`$ is the wave front curvature radius. This parameter fully characterizes the beam and satisfies the Kogelnik $`ABCD`$ law; namely, after propagation through a first-order system, the parameter $`Q`$ changes to $`Q^{}`$ via
$$Q^{}=\mathrm{\Psi }[𝐌,Q]=\frac{C+DQ}{A+BQ}.$$
(2.8)
Since $`ImQ>0`$ by the definition (2.6), one immediately checks that $`ImQ^{}>0`$ and we can thus view the action of the first-order system as a bilinear transformation $`\mathrm{\Psi }`$ on the upper complex half-plane. When we use the metric $`ds=|dQ|/ImQ`$ to measure distances, what we get is the standard model of the hyperbolic plane $``$ . This plane $``$ is invariant under bilinear transformations.
We note that the whole real axis, which is the boundary of $``$, is also invariant under (2.8) and represents wave fields with unlimited transverse irradiance (contrary to the notion of a beam). On the other hand, for the points in the imaginary axis we have an infinite wave front radius, which defines the corresponding beam waists. The origin represents a plane wave.
Bilinear transformations constitute an important tool in many branches of physics. For example, in polarization optics they have been employed for a simple classification of polarizing devices by means of the concept of eigenpolarizations of the transfer function .
In our context, the equivalent concept can be stated as the beam configurations such that $`Q=Q^{}`$ in equation (2.8), whose solutions are
$$Q_\pm =\frac{1}{2B}\left[(DA)\pm \sqrt{(A+D)^24}\right].$$
(2.9)
These values of $`Q`$ are known as the fixed points of the transformation.
The trace of $`𝐌`$, $`\mathrm{Tr}(𝐌)=A+D`$, provides a suitable tool for the classification of optical systems . It has also played an important role in studying propagation in periodic media . When $`[\mathrm{Tr}(𝐌)]^2<4`$ the action is said elliptic and there are no real roots: they are complex conjugates and only one of them lies in $``$, while the other lies outside. When $`[\mathrm{Tr}(𝐌)]^2>4`$ there are two real roots (i.e., in the boundary of $``$) and the action is hyperbolic. Finally, when $`[\mathrm{Tr}(𝐌)]^2=4`$ there is one (double) real solution and the system action is called parabolic.
To proceed further let us note that by taking the conjugate of $`𝐌`$ with any matrix $`𝐂`$ SL($`2,)`$
$$𝐌_\mathrm{C}=𝐂𝐌𝐂^1,$$
(2.10)
we obtain another matrix of the same type, since $`\mathrm{Tr}(𝐌)=\mathrm{Tr}(𝐌_\mathrm{C})`$. Conversely, if two systems have the same trace, one can always find a matrix $`𝐂`$ satisfying equation (2.10).
Note that $`Q`$ is a fixed point of $`𝐌`$ if and only if the image of $`Q`$ by $`𝐂`$ (i.e., $`\mathrm{\Psi }[𝐂,Q]`$) is a fixed point of $`𝐌_\mathrm{C}`$. In consequence, given any ray-transfer matrix $`𝐌`$ one can always find a $`𝐂`$ such that $`𝐌_\mathrm{C}`$ takes one of the following canonical forms :
$`𝐊_\mathrm{C}(\vartheta )`$ $`=`$ $`\left(\begin{array}{cc}\mathrm{cos}(\vartheta /2)& \mathrm{sin}(\vartheta /2)\\ \mathrm{sin}(\vartheta /2)& \mathrm{cos}(\vartheta /2)\end{array}\right),`$ (2.13)
$`𝐀_\mathrm{C}(\xi )`$ $`=`$ $`\left(\begin{array}{cc}e^{\xi /2}& 0\\ 0& e^{\xi /2}\end{array}\right),`$ (2.16)
$`𝐍_\mathrm{C}(\nu )`$ $`=`$ $`\left(\begin{array}{cc}1& 0\\ \nu & 1\end{array}\right),`$ (2.19)
where $`0\vartheta 4\pi `$ and $`\xi ,\nu `$. These matrices define the one-parameter subgroups of SL(2, $``$) and have as fixed points $`+i`$ (elliptic), 0 and $`\mathrm{}`$ (hyperbolic), and $`\mathrm{}`$ (parabolic), respectively. They are the three basic blocks in terms of which any system action can be expressed. Clearly, $`𝐊_\mathrm{C}(\vartheta )`$ represents a rotation in phase space, $`𝐀_\mathrm{C}(\xi )`$ is a magnifier that scales $`x`$ up by the factor $`e^{\xi /2}`$ and $`p`$ down by the same factor, and $`𝐍_\mathrm{C}(\nu )`$ represents the action of a thin lens of power $`\nu `$ (i.e., focal length $`1/\nu `$.
For the canonical forms (2.13), the corresponding actions are
$`Q^{}`$ $`=`$ $`{\displaystyle \frac{\mathrm{cos}(\vartheta /2)Q\mathrm{sin}(\vartheta /2)}{\mathrm{sin}(\vartheta /2)Q+\mathrm{cos}(\vartheta /2)}},`$
$`Q^{}`$ $`=`$ $`e^\xi Q,`$ (2.20)
$`Q^{}`$ $`=`$ $`Q+\nu .`$
The first is a rotation, in agreement with Euclidean geometry, since a rotation has only one invariant point. The second is a translation because it has no fixed points in $``$ and the geodesic line joining the two fixed points (0 and $`\mathrm{}`$) remains invariant (it is the axis of the translation). The third one is known as a parallel displacement.
When one of the parameters $`\theta `$, $`\xi `$, or $`\nu `$ in (2) varies, the transformed points $`Q^{}`$ describe a curve called the orbit of $`Q`$ under the action of the corresponding one-parameter subgroup. In figure 1.a we have plotted typical orbits for the canonical forms (2.13). For matrices $`𝐊_\mathrm{C}(\vartheta )`$ the orbits are circumferences centered at the invariant point $`+i`$ and passing through $`Q`$ and $`1/Q`$. For $`𝐀_\mathrm{C}(\xi )`$, they are lines going from 0 to the $`\mathrm{}`$ through $`Q`$ and they are known as hypercicles. Finally, for matrices $`𝐍_\mathrm{C}(\nu )`$ the orbits are lines parallel to the real axis passing through $`Q`$ and they are known as horocycles .
For a general matrix $`𝐌`$ the corresponding orbits can be obtained by transforming with the appropriate matrix $`𝐂`$ the orbits described before. The explicit construction of the family of matrices $`𝐂`$ is not difficult: it suffices to impose that $`𝐂`$ transforms the fixed points of $`𝐌`$ into the ones of $`𝐊_\mathrm{C}(\vartheta )`$, $`𝐀_\mathrm{C}(\xi )`$, or $`𝐍_\mathrm{C}(\nu )`$, respectively. Just to work out an example that will play a relevant role in the forthcoming, we consider a matrix $`𝐌`$ representing an elliptic action with one fixed point denoted by $`Q_f`$. Since the fixed point for the corresponding canonical matrix $`𝐊_\mathrm{C}(\vartheta )`$ is $`+i`$, the matrix $`𝐂`$ we are looking for is determined by
$$\mathrm{\Psi }[𝐂,Q_f]=i.$$
(2.21)
If the matrix $`𝐂`$ is written as
$$𝐂=\left(\begin{array}{cc}C_1& C_2\\ C_3& C_4\end{array}\right),$$
(2.22)
the solution of (2.21) is
$`C_2`$ $`=`$ $`{\displaystyle \frac{C_1ReQ_f+C_3ImQ_f}{|Q_f|^2}},`$
$`C_4`$ $`=`$ $`{\displaystyle \frac{C_1ImQ_fC_3ReQ_f}{|Q_f|^2}}.`$
In addition, the condition $`det𝐂=+1`$ imposes
$$C_3=\sqrt{\frac{|Q_f|^2}{ImQ_f}C_1^2},$$
(2.24)
that, together (2) determines the matrix $`𝐂`$ in terms of the free parameter $`C_1`$.
In figure 1.b we have plotted typical examples of such orbits for elliptic, hyperbolic, and parabolic actions. We stress that once the fixed points of the ray-transfer matrix are known, one can ensure that $`Q^{}`$ will lie in the orbit associated to $`Q`$.
## 3 First-order systems as transformations in the Poincaré unit disc $`𝔻`$
To complete the geometrical setting introduced in the previous Section, we explore now a remarkable transformation (introduced by Cayley) that maps bijectively the hyperbolic plane $``$ onto the unit disc, denoted by $`𝔻`$. This can be done via the unitary matrix
$$𝓤=\frac{1}{\sqrt{2}}\left(\begin{array}{cc}1& i\\ i& 1\end{array}\right),$$
(3.1)
in such a way that
$$𝓜=𝓤𝐌𝓤^1=\left(\begin{array}{cc}\alpha & \beta \\ \beta ^{}& \alpha ^{}\end{array}\right),$$
(3.2)
where $`𝓜`$ is a matrix with $`det𝓜=+1`$ and whose elements are given in terms of those of $`𝐌`$ by
$`\alpha `$ $`=`$ $`{\displaystyle \frac{1}{2}}[(A+D)+i(CB)],`$
$`\beta `$ $`=`$ $`{\displaystyle \frac{1}{2}}[(B+C)+i(DA)].`$
In other words, the matrices $`𝓜`$ belong to the group SU(1, 1), which plays an essential role in a variety of branches in physics. Obviously, the bilinear action induced by these matrices is
$$𝒬^{}=\mathrm{\Phi }[𝓜,𝒬]=\frac{\beta ^{}+\alpha ^{}𝒬}{\alpha +\beta 𝒬},$$
(3.4)
where $`𝒬`$ is the point transformed by (3.1) of the original $`Q`$:
$$𝒬=\frac{Qi}{1iQ}.$$
(3.5)
The transformation by $`𝓤`$ establishes then a one-to-one map between the group SL(2, $``$) of matrices $`𝐌`$ and the group SU(1, 1) of complex matrices $`𝓜`$, which allows for a direct translation of the properties from one to the other.
It is easy to see that $``$ maps onto $`𝔻`$, as desired. The imaginary axis in $``$ goes to the $`Y`$ axis of the disc $`𝔻`$ (in both cases, $`R=\mathrm{}`$ and define beam waists). In particular, $`Q=+i`$ is mapped onto $`𝒬=0`$. The boundary of $``$ (the real axis) goes to the boundary of $`𝔻`$ (the unit circle), and both boundaries represent fully unlimited irradiance distributions (i.e., non-beam solutions).
Since the matrix conjugation (3.2) does not change the trace, the same geometrical classification in three basic actions still holds. In fact, by conjugating with $`𝓤`$ the canonical forms (2.13), we get the corresponding ones for SU(1, 1):
$`𝓚_𝒞(\vartheta )`$ $`=`$ $`\left(\begin{array}{cc}\mathrm{exp}(i\vartheta /2)& 0\\ 0& \mathrm{exp}(i\vartheta /2)\end{array}\right),`$ (3.8)
$`𝓐_𝒞(\xi )`$ $`=`$ $`\left(\begin{array}{cc}\mathrm{cosh}(\xi /2)& i\mathrm{sinh}(\xi /2)\\ i\mathrm{sinh}(\xi /2)& \mathrm{cosh}(\xi /2)\end{array}\right),`$ (3.11)
$`𝓝_𝒞(\nu )`$ $`=`$ $`\left(\begin{array}{cc}1i\nu /2& \nu /2\\ \nu /2& 1+i\nu /2\end{array}\right),`$ (3.14)
that have as fixed points 0 (elliptic), $`+i`$ and $`i`$ (hyperbolic) and $`+i`$ (parabolic), respectively. The first matrix represent a rotation in phase space, also called a fractional Fourier transformation, while the second one is sometimes called a hyperbolic expander .
The corresponding orbits for these matrices are defined by
$`𝒬^{}`$ $`=`$ $`\mathrm{\Phi }[𝓚_𝒞,𝒬]=𝒬\mathrm{exp}(i\vartheta ),`$
$`𝒬^{}`$ $`=`$ $`\mathrm{\Phi }[𝓐_𝒞,𝒬]={\displaystyle \frac{𝒬i\mathrm{tanh}(\xi /2)}{1+i𝒬\mathrm{tanh}(\xi /2)}},`$ (3.15)
$`𝒬^{}`$ $`=`$ $`\mathrm{\Phi }[𝓝_𝒞,𝒬]={\displaystyle \frac{𝒬+(1+i𝒬)\nu /2}{1+(𝒬i)\nu /2}}.`$
As plotted in figure 2.a, for matrices $`𝓚_𝒞(\vartheta )`$ the orbits are circumferences centered at the origin. For $`𝓐_𝒞(\xi ),`$ they are arcs of circumference going from the point $`+i`$ to the point $`i`$ through $`𝒬`$. Finally, for the matrices $`𝓝_𝒞(\nu )`$ the orbits are circumferences passing through the points $`i`$, $`𝒬`$, and $`𝒬^{}`$. In figure 2.b we have plotted the corresponding orbits for arbitrary fixed points.
## 4 Application to optical resonators
The geometrical ideas presented before allows one to describe the evolution of a GSM beam by means of the associated orbits. As an application of the formalism, we consider the illustrative example of an optical cavity consisting of two spherical mirrors of radii $`R_1`$ and $`R_2`$, separated a distance $`d`$. The ray-transfer matrix corresponding to a round trip can be routinely computed
$$𝐌=\left(\begin{array}{cc}2g_1g_2g_1+g_21& \frac{d}{2}(2g_1g_2+g_1+g_2)\\ \frac{2}{d}(2g_1g_2g_1g_2)& 2g_1g_2+g_1g_21\end{array}\right),$$
(4.1)
where we have used the parameters ($`i=1,2`$)
$$g_i=1\frac{d}{R_i}.$$
(4.2)
Note that
$$\mathrm{Tr}(𝐌)=2(2g_1g_21).$$
(4.3)
Since the trace determines the fixed point and the orbits of the system, the $`g`$ parameters establish uniquely the geometrical action of the resonator. To clarify further this point, in figure 3 we have plotted the value of $`|\mathrm{Tr}(𝐌)|`$ in terms of $`g_1`$ and $`g_2`$. The plane $`|\mathrm{Tr}(𝐌)|=2`$, which determines the boundary between elliptic and hyperbolic action, is also shown. At the top of the figure, a density plot is presented, with the characteristic hyperbolic contours.
Assume now that the light bounces $`N`$ times through this system. The overall transfer matrix is then $`𝐌^N`$, so all the algebraic task reduces to finding a closed expression for the $`N`$th power of the matrix $`𝐌`$. Although there are several elegant ways of computing this power , we shall instead apply our geometrical picture: the transformed beam is represented by the point
$$Q_N=\mathrm{\Psi }[𝐌,Q_{N1}]=\mathrm{\Psi }[𝐌^N,Q_0],$$
(4.4)
where $`Q_0`$ denotes the initial point.
Note that all the points $`Q_N`$ lie in the orbit associated to the initial point $`Q_0`$ by the single round trip, which is determined by its fixed points: the character of these fixed points determine thus the behaviour of this periodic system. By varying the parameters $`g`$ of the resonator we can choose to work in the elliptic, the hyperbolic, or the parabolic case .
To illustrate how this geometrical approach works in practice, in figure 4.a we have plotted the sequence of successive iterates obtained for different kind of ray-transfer matrices, according to our previous classification. In figure 4.b we have plotted the same sequence but in the unit disc, obtained via the unitary matrix $`𝓤`$.
In the elliptic case, it is clear that the points $`Q_N`$ revolve in the orbit centered at the fixed point and the system never reaches the real axis. Equivalently, the points $`𝒬_N`$ never reach the unit circle.
On the contrary, for the hyperbolic and parabolic cases the iterates converge to one of the fixed points on the real axis, although with different laws . In the general context of scattering by periodic systems this corresponds to the band stop and band edges, respectively .
What we conclude from this analysis is that the iterates of hyperbolic and parabolic actions produce solutions fully unlimited, which are incompatible with our ideas of a beam. The only beam solutions are thus generated by elliptic actions and, according with equation (4.3), the stability criterion is
$$0|2g_1g_21|=|\mathrm{cos}(\vartheta /2)|1,$$
(4.5)
where $`\vartheta `$ is the parameter in the canonical form $`𝐊_\mathrm{C}`$ in equation (2.13). Such a condition is usually worked out in terms of algebraic arguments using ray-transfer matrices, although the final results apply exclusively to scalar wave fields.
Finally, we stress that real cavities resonate with vector fields. The situation then is far more involved because the vector diffraction for (polarized) electric fields is more difficult to handle, even for systems with small Fresnel numbers and the $`ABCD`$ law does not apply to the corresponding kernel . Exact solutions for these vector beams have recently appeared . In any case, there is abundant evidence that the stability condition (4.5) works well. This could be expected, since the transition to scalar theories captures all the essential physics embodied in the more elaborated vector analogues .
## 5 Concluding remarks
In this paper, we have provided a geometrical scenario to deal with first-order optical systems. More specifically, we have reduced the action of any system to a rotation, a translation or a parallel displacement, according to the magnitude of the trace of its ray-transfer matrix. These are the basic isometries of the hyperbolic plane $``$ and also of the Poincaré unit disc $`𝔻`$. We have also provided an approach for a qualitative examination of the stability condition of an optical resonator.
We hope that this approach will complement the more standard algebraic techniques and together they will help to obtain a better physical and geometrical feeling for the properties of first-order optical systems. |
warning/0506/astro-ph0506740.html | ar5iv | text | # Mapping dark matter with cosmic magnification
## Abstract
We develop a new tool to generate statistically precise dark matter maps from the cosmic magnification of galaxies with distance estimates. We show how to overcome the intrinsic clustering problem using the slope of the luminosity function, because magnificability changes strongly over the luminosity function, while intrinsic clustering only changes weakly. This may allow precision cosmology beyond most current systematic limitations. SKA is able to reconstruct projected matter density map at smoothing scale $`10^{^{}}`$ with S/N$`1`$, at the rate of $`200`$-$`4000`$ deg<sup>2</sup> per year, depending on the abundance and evolution of 21cm emitting galaxies. This power of mapping dark matter is comparable to, or even better than that of cosmic shear from deep optical surveys or 21cm surveys.
preprint: FERMILAB-PUB-05-260-A
Introduction.— The precision mapping of the universe, and the accurate determination of cosmological parameters have been enabled by the recent generation of cosmic microwave background(CMB) experiments, galaxy and lensing surveys, and new analysis techniques. Weak gravitational lensing has emerged with a promising future of mapping dark matter directly, which would allow the inference of the state of the universe, including its dynamics and the nature of dark energy. Lensing is free from modeling assumptions, and can be accurately predicted from first principles. Several major surveys are underway, under construction or in the planning stage. Currently, most attention has focused on using the lensing induced cosmic shear. But such an approach is subject to a series of difficult experimental systematics . CMB lensing and 21cm background lensing are promising. But contaminations such as the kinetic Sunyaev Zeldovich effect and/or non-Gaussianity may degrade their accuracy. In this paper we will address an alternative approach, the lensing induced cosmic magnification, which is not subject to the known problems, and could provide a robust statistical signal.
Traditionally, intrinsic clustering had presented a serious problem to measurement of cosmic magnification. The observable quantity is the surface density of galaxies above some flux threshold. A variation in this surface density is then interpreted as lensing. Unfortunately, intrinsic clustering is usually larger than the lensing induced signal. By utilizing the redshift information, intrinsic clustering can be effectively eliminated in lensing correlation functions. In this paper, we further show that, beyond the above statistical lensing measurement, 2D convergence $`\kappa `$ maps can be reconstructed with lower systematics and larger sky coverage than cosmic shear maps, by utilizing both the redshift and flux information of galaxies. 2D $`\kappa `$ maps not only provide independent and robust constraints on cosmology, but also are complementary to traditional shear maps. It allows one to explicitly and locally solve for non-reduced shear, an independent mode of checking E-B decomposition, and break the mass-sheet degeneracy.
Cosmic magnification.— Cosmic magnification causes coherent changes in the apparent galaxy number density. Let $`N_{ij}`$ be the observed number of galaxies (including false peaks) at the $`i`$-th flux bin and $`j`$-th redshift bin, falling into an angular pixel centered at direction $`\widehat{n}`$ with angular size $`\theta `$. It can be expressed as
$`N_{ij}(\widehat{n})=\overline{N}_{ij}+\overline{N}_{ij}^r\left[W_{ij}\kappa _j(\widehat{n})+\delta _{g,ij}(\widehat{n})\right]+\delta N_{P,ij}(\widehat{n}).`$ (1)
The signal $`W\kappa `$ has unique dependence on galaxy flux through $`W=2(\alpha 1)`$. Here, $`\alpha =d\mathrm{ln}[dn/dF]/d\mathrm{ln}F1`$ and $`dn/dF`$ is the mean number of observed galaxies per flux interval <sup>1</sup><sup>1</sup>1The observed $`dn/dF`$ is convolved with system noise. Because there are more dwarf galaxies than massive ones, noise makes the observed $`dn/dF`$ both larger and steeper, in the flux range that SKA can probe at $`z2`$. The overall effect is that system noise in flux measurements increases the cosmic magnification signal and strengthens the result in this paper. For simplicity, we neglect this complexity. $`\overline{N}_{ij}=\overline{N}_{ij}^r+\overline{N}_{ij}^f`$, $`\overline{N}_{ij}^r`$, $`\overline{N}_{ij}^f`$ are the mean number of detections, real galaxies and false peaks, respectively. $`\delta _g`$ and $`\delta N_P`$ are galaxy intrinsic clustering and Poisson fluctuation, respectively.
Our goal is to recover $`\kappa `$ of each angular pixel, given observables $`N_{ij}`$, $`\overline{N}_{ij}^r`$, $`W_{ij}`$ and $`\overline{N}_{ij}`$ <sup>2</sup><sup>2</sup>2Cosmic magnification does not change the averaged galaxy spatial and flux distribution, up to $`O(\kappa ^2)10^4`$ accuracy. The sky coverage of SKA is $`100`$ deg<sup>2</sup>. Thus for each redshift and flux bin, there are $`4000`$ angular pixels with size $`\theta 10^{^{}}`$ and $`10^5`$ galaxies across the survey sky, so $`\overline{N}_{ij}`$ can be measured accurately. $`\overline{N}_{ij}^f`$ can be accurately predicted, since system noise is Gaussian and the dispersion $`S_{\mathrm{sys}}`$ is specified for each survey. The number of false peaks with flux above $`n`$-$`\sigma `$, or $`nS_{\mathrm{sys}}`$ per redshift interval per beam is $`[1.4\mathrm{Ghz}/\mathrm{\Delta }\nu (1+z)^2]\mathrm{Erfc}[n/\sqrt{2}]/2`$. $`\mathrm{\Delta }\nu `$ is chosen to be the frequency width corresponding to $`100`$ km/$`s`$ velocity dispersion at redshift z. Thus, one can accurately predict $`\overline{N}_{ij}^r`$ and $`W_{ij}`$.. We consider SKA<sup>3</sup><sup>3</sup>3SKA:http://www.skatelescope.org/, which can detect $`10^8`$ high z galaxies through the neutral hydrogen 21cm emission line. $`\kappa `$ has typical value $`1\%`$. To beat down Poisson fluctuations, $`10^4`$ galaxies per angular pixel are required. Traditionally, objects are selected at a 5$`\sigma `$ cut, where one can neglect the fraction of false detections. This of course also discards the majority of the signal. With a $`0.5`$-$`\sigma `$ cut, one can reduce Poisson noise at $`\theta 10^{^{}}`$. To increase lensing signal while reducing $`\delta _g`$ contamination, we focus on source redshifts $`z2`$. After averaging over the full redshift range $`z2`$, $`\delta _g`$ is still several times larger than $`\kappa `$. However, $`\delta N_P`$ and $`\delta _{g,ij}`$ have different flux dependence to that of the signal. Weighting each galaxies by some function of their flux can suppress the prefactors of $`\delta _g`$ and $`\delta N_P`$. Intuitively, Eq. 1 implies the optimal estimator to be linear in $`N_{ij}`$.
The predictions rely on the assumed HI mass function $`n(M_{\mathrm{HI}},z)`$. We extrapolate the locally observed $`n(M_{\mathrm{HI}},z)=n_0(z)(M_{\mathrm{HI}}/M_{})^{1.2}\mathrm{exp}(M_{\mathrm{HI}}/M_{})`$ to high redshifts either assuming no evolution in both $`n_0`$ and $`M_{}`$ (conservative case) or $`n_0(z),M_{}(z)(1+z)^{1.45}\mathrm{exp}(z/2.6)`$ (realistic case), which is calibrated against Lyman-$`\alpha `$ observations (refer to for details). We adopt a flat $`\mathrm{\Lambda }`$CDM cosmology with $`\mathrm{\Omega }_m=0.3`$, $`\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$, $`h=0.7`$, $`\sigma _8=0.9`$, the primordial power index $`n=1`$, BBKS transfer function and Peacock-Dodds fitting formula for the nonlinear density power spectrum .
The optimal estimator.— Since $`z2`$ galaxies are mainly lensed by matter at $`z1`$, $`\kappa =AB/\chi (z)`$ is an excellent approximation, where $`A`$ and $`B`$ are two constants and $`\chi `$ is the comoving angular diameter distance. Since $`\chi (z)`$ varies slowly at $`z>2`$, one can approximate $`\kappa (\chi ,\widehat{n})\kappa =\kappa (\chi ,\widehat{n})`$, where $`\chi =\overline{N}_{ij}^r/_{ij}\chi _j^1\overline{N}_{ij}^r`$ is the effective distance to lens <sup>4</sup><sup>4</sup>4The approximation $`\kappa \kappa `$ simplifies the derivation of the optimal estimator significantly, though its accuracy can be as bad as $`\pm 20\%`$, at each redshift bins. But after averaging over many redshift bins, corrections in different bins effectively cancel. For the optimal estimators derived (Eq. 3 & 4), one can derive the unbiased expression of $`\kappa `$ such that $`\kappa =\widehat{\kappa }`$.. In the limit that $`\overline{N}_{ij}1`$, Poisson fluctuations become Gaussian. The likelihood function of $`\kappa `$ at an angular pixel, marginalized over $`p(\delta _{g,11}\mathrm{}\delta _{g,ij})`$, the probability distribution of $`\delta _{g,ij}`$ of this angular pixel, is
$`L`$ $``$ $`{\displaystyle \mathrm{exp}\left[\underset{ij}{}\frac{[N_{ij}\overline{N}_{ij}\overline{N}_{ij}^r(W_{ij}\kappa +\delta _{g,ij})]^2}{2\overline{N}_{ij}}\right]}`$ (2)
$`\times p(\delta _{g,11}\mathrm{},\delta _{g,ij}){\displaystyle \underset{ij}{}}d\delta _{g,ij}.`$
We choose the redshift bin size $`\mathrm{\Delta }z0.2`$ and angular pixel size $`10^{^{}}`$ such that $`\delta _{g,ij}`$ of different redshift bins are uncorrelated. For this choice, the matter density dispersion of each redshift bin $`\sigma _m0.1`$. This verifies the neglect of high order term $`\delta _g\kappa `$ in Eq. 1 & 2.
Since $`\sigma _m0.1`$ and galaxy bias $`b_g`$ is unlikely bigger than several, it is reasonable to assume that galaxies are Gaussian distributed. Then $`p(\delta _{g,11}\mathrm{})`$ is completely determined by the covariance matrix $`C_{i_1j_1;i_2j_2}\delta _{g,i_1j_1}(\widehat{n})\delta _{g,i_2j_2}(\widehat{n})`$. SKA can directly and accurately measure the correlations of galaxy density fluctuations between flux bins, which are the sum of $`C_{i_1j_1;i_2j_2}`$, correlations induced by lensing and cross terms. In the interesting range, $`C_{i_1j_1i_2j_2}`$ dominates. So one can take the measured sum as first guess of $`C_{i_1j_1;i_2j_2}`$. Maximizing $`L`$, one obtains the optimal estimator $`\widehat{\kappa }`$ of $`\kappa `$. The reconstructed $`\kappa `$ can in turn be applied to subtract the lensing contribution in the covariance matrix estimation. This can be done iteratively. Since the lensing contribution is small, such iteration should be stable and converge quickly.
The properties of high redshift 21cm emitting galaxies are currently poorly known. It is likely that they trace the underlying dark matter at some level, and that galaxies of different luminosities are correlated to each other. We consider this case first, and then the extreme stochastic biasing limit where galaxies of different flux are uncorrelated with each other. These two cases correspond to the worst and best cases for the $`\kappa `$ reconstruction, respectively.
Deterministic biasing.— We first consider the case that $`\delta _{g,ij}`$ of different flux bins (but of the same redshift bin) are linearly correlated, namely, $`\delta _{g,ij}=b_{ij}\delta _j`$, where $`\delta _j`$ is the dark matter density of the $`j`$-th redshift bin. As discussed above, $`b_{ij}`$ can be measured iteratively. Marginalizing over $`\delta _j`$, we obtain
$`L`$ $``$ $`\mathrm{exp}\left[{\displaystyle \frac{(\kappa \widehat{\kappa })^2}{2(\mathrm{\Delta }\kappa )^2}}\right],`$
$`\widehat{\kappa }`$ $`=`$ $`({\displaystyle \underset{j}{}}S_j{\displaystyle \frac{B_jQ_j}{A_j}})(\mathrm{\Delta }\kappa )^2,`$
$`\mathrm{\Delta }\kappa `$ $`=`$ $`\left({\displaystyle \underset{j}{}}T_j{\displaystyle \frac{Q_j^2}{A_j}}\right)^{1/2}`$
$``$ $`\left(\overline{N}W^2\overline{N}{\displaystyle \frac{Wb^2}{b^2}}\right)^{1/2}.`$
Here $`A_j=_i(\overline{N}_{ij}^rb_{ij})^2/\overline{N}_{ij}+1/\sigma _j^2`$, $`B_j=_i(N_{ij}\overline{N}_{ij})\overline{N}_{ij}^rb_{ij}/\overline{N}_{ij}`$, $`Q_j=_i\overline{N}_{ij}^{r,2}W_{ij}b_{ij}/\overline{N}_{ij}`$, $`S_j=_i(N_{ij}\overline{N}_{ij})\overline{N}_{ij}^rW_{ij}/\overline{N}_{ij}`$, and $`T_j=_i(\overline{N}_{ij}^rW_{ij})^2/\overline{N}_{ij}`$. $`\overline{N}`$ is the mean number of galaxies in each angular pixel. $`\mathrm{}`$ are weighted by galaxies with the noise from false peaks taken into account.
Maximal stochasticity.— Stochasticity eases the subtraction of the intrinsic clustering signal. In this case, $`\delta _g`$ of different bins are uncorrelated. We have
$`L`$ $``$ $`\mathrm{exp}\left[{\displaystyle \underset{ij}{}}{\displaystyle \frac{[N_{ij}\overline{N}_{ij}\overline{N}_{ij}^rW_{ij}\kappa ]^2}{2\sigma _{ij}^2}}\right],`$
$`\widehat{\kappa }`$ $`=`$ $`{\displaystyle \frac{_{ij}(N_{ij}\overline{N}_{ij})\overline{N}_{ij}^rW_{ij}/\sigma _{ij}^2}{_{ij}[\overline{N}_{ij}^rW_{ij}]^2/\sigma _{ij}^2}},`$
$`\mathrm{\Delta }\kappa `$ $`=`$ $`[^2\mathrm{ln}L/\kappa ^2]^{1/2}=\left[{\displaystyle \underset{ij}{}}{\displaystyle \frac{[\overline{N}_{ij}^rW_{ij}]^2}{\sigma _{ij}^2}}\right]^{1/2}`$ (4)
$``$ $`[\overline{N}W^2]^{1/2}{}_{}{}^{}=_{}^{}when\overline{N}_{ij}^r\sigma _{g,ij}^20.`$
Here, $`\sigma _{ij}^2=\overline{N}_{ij}+\overline{N}_{ij}^{r,2}\sigma _{g,ij}^2`$, where the first term is the shot noise and the second term is the intrinsic fluctuation of galaxy number distribution. The conditions $`\overline{N}_{ij}^r\sigma _{g,ij}^20.01\overline{N}_{ij}0`$ and $`\overline{N}_{ij}1`$ (for Gaussianity) can both be satisfied since galaxy bias $`b_g`$ is unlikely bigger than several. A similar estimator has been derived by . In two estimators, $`W^2`$ and $`Wb`$ are two key ingredients and reflect the key role of flux information.
Results.— SKA is able to detect $`n_g100\mathrm{arcmin}^2`$ galaxies at $`z2`$ (table I). For an integration time $`t_{\mathrm{int}}=18`$ days/deg<sup>2</sup>, a S/N$`2`$ can be achieved at $`\theta 10^{^{}}`$ (fig. 1). Deep survey configuration detects more faint galaxies, which have $`W2`$, mimic a constant $`b`$ and thus do not contribute to the signal, due to the $`W^2Wb^2/b^2`$ facotr in Eq. 3. An optimal survey configuration should have $`Wb0`$, which can be achieved at $`t_{\mathrm{int}}0.2`$-$`1`$ day/deg<sup>2</sup> (fig. 1). Since $`n_g`$ above $`0.5`$-$`\sigma `$ decreases much more slowly than $`t_{\mathrm{int}}`$ (for example, for the evolution model, decreasing $`t_{\mathrm{int}}`$ from $`180`$ days/deg<sup>2</sup> to $`4`$ hours/deg<sup>2</sup>, $`n_g`$ only decreases by a factor of $`9`$), it is still likely to achieve S/N$`>1`$ at $`\theta 10^{^{}}`$ and scan rate of $`4000`$ deg<sup>2</sup> per year (fig. 2). This will produce more lensing information ($``$ S/N$`\times f_{\mathrm{sky}}^{1/2}`$) in a one year SKA survey than SNAP<sup>5</sup><sup>5</sup>5SNAP: http://snap.lbl.gov/ will produce, which will cover $`1000`$ deg<sup>2</sup> sky area with S/N$`2`$ at smoothing scale $`\theta 10^{^{}}`$.
Since the SKA will have $`0.3^{^{\prime \prime }}`$ resolution at $`z2`$, it can resolve galaxies and measure cosmic shear. An intrinsic advantage of cosmic magnification measurement over cosmic shear measurement is that it does not require galaxies to be resolved. Thus, dwarf galaxies which are too small and too faint for reliable shear measurement still contribute to magnification measurement. Cosmic magnification exceeds cosmic shear at integration rate $`0.2`$-$`10`$ days/deg<sup>2</sup> (fig.1). We note that this comparison is conservative. We have neglected all systematics of shear measurement. For magnification estimation, we only select galaxies above a $`0.5`$-$`\sigma `$ detection threshold, or HI mass above $`\mathrm{several}\times 10^8M_{}h^2`$. There are numerous galaxies with HI mass $`10^7M_{}h^2`$, which can in principle be used to improve the measurement. We do not explore its potential in this paper since the luminosity function at the faint end is unclear.
Several uncertainties could degrade the signal separation. (1) The HI mass function, which is the dominant factor, as can be seen from table I and fig. 1. Here we further draw the attention on the slope of the HI mass function. For an extreme case that $`\alpha 1`$ and $`W0`$ over a large flux range, the signal disappears. This effect can be straightforwardly estimated through the $`W^2`$ and $`Wb`$ terms in Eq. 3 & 4, once the HI mass function is measured. Since HI mass function at high $`z`$ is effectively unknown, we postpone the discussion in this paper. (2) The galaxy bias. For the case of deterministic biasing, if $`b_gW`$, flux information is no longer useful for the separation and our method effectively fails. But since $`b_g>0`$, as long as the survey is deep enough to probe the faint end of galaxies where $`W<0`$, $`b_g`$ can not always mimic $`W`$ and the separation is always possible. (3) The galaxy distribution. When $`b_g`$ is bigger than several or smoothing size is smaller than several arc-minutes, $`\delta _g`$ is non-Gaussian. In this case, the estimators described above are no longer optimal. Optimal estimators for non-Gaussian galaxy distribution should be further investigated.
Applications.—The reconstructed $`\kappa `$ map can be applied to measure many lensing statistics. For this purpose, reconstructed $`\kappa `$ can be noisy because these statistics generally average over many angular pixels and achieve high S/N. Then the optimal estimator derived in this paper can be applied to each narrow redshift bins and allows the lensing tomography. (1)The probability density function $`p(\kappa )`$. $`p(\kappa )`$ as a function of $`\kappa `$ and smoothing angular size $`\theta `$ can provide independent constraints on cosmology. Recently showed that the Wiener filter reconstruction of $`p(\kappa )`$ from noisy convergence map can go deep into regions where $`|\kappa /\mathrm{\Delta }\kappa |1`$. We thus expect that $`p(\kappa )`$ can be recovered accurately from SKA. (2) Lensing power spectrum and bispectrum. The reconstructed $`\kappa `$ map barely has S/N$`5`$, so it is consistent to neglect $`10\%`$ higher order terms: $`O(\kappa ^2)`$ terms and $`\delta _g\kappa `$ term neglected in Eq. 1 and $`\kappa (\chi )\kappa `$. But these terms should be taken into account for precision measurement of lensing power spectrum and bispectrum, since their statistical errors can reach $`1\%`$ accuracy. For the linear estimator we derived, contributions of these terms to the power spectrum and bispectrum can be straightforwardly and robustly predicted. So, there is no need to derive a more complicated nonlinear estimator. (3) Cluster finding and cluster density profile. This is a promising approach to break the cluster mass sheet degeneracy. In the reconstructed maps, massive clusters at $`z0.2`$ show as high peaks with strength $`\kappa 0.1`$ and size $`10^{^{}}`$ and can be easily identified. These clusters are excellent objects to measure the geometry of the universe by the technique of lensing cross-correlation tomography . Since S/N is so high, one can choose smoothing size $`1^{^{}}`$ and measure the projected cluster density profile. Exerting a prior on cluster density profile, the reconstruction can be further improved . When $`\kappa 1`$, the weak lensing condition breaks and Eq. 1 no longer holds. By utilizing the exact magnification equation, one can develop new estimator, in analogy to the reduced shear reconstruction . We leave this topic for further study.
We summarize our results. Cosmic magnification is statistically more sensitive than cosmic shear because it is possible to use the large number of galaxies detected at low statistical significance. Intrinsic clustering can be subtracted because (1) magnification depends strongly on the shape of the luminosity function, which varies significantly, while intrinsic clustering depends weakly on the intrinsic luminosity itself and (2) they have different redshift dependence. Cosmic magnification shows promise as a complementary technique to map the statistically precise distribution of matter, which is not subject to most of the systematics of cosmic shear. We have worked through the specific numbers for the SKA, but the general formalism would also apply to optical spectroscopic or photometric redshift surveys.
Acknowledgments.— We thank Scott Dodelson for many helpful conversations and careful proofreading. We thank Albert Stebbins and Martin White for helpful discussions. P.J. Zhang was supported by the DOE and the NASA grant NAG 5-10842 at Fermilab. |
warning/0506/quant-ph0506231.html | ar5iv | text | # Einstein’s Photon Concept Quantified by the Bohr Model of the Photon
## 1 Einstein’s Concept of the Photon
In 1905 Einstein published a celebrated paper popularly known as “the photo-electric paper”, for which he was awarded the Nobel prize some 16 years later . While this paper explains other physical phenomena in addition to the photo-electric effect, the unifying concept is that:
> “The energy in a beam of light is not uniformly distributed as in a classical plane wave, but is localized in packets of electromagnetic radiation, each packet having an energy $`h\nu =hc/\lambda `$, where $`h`$ is Planck’s constant, and $`\nu `$ and $`\lambda `$ are the frequency and wavelength of the radiation.”
Einstein called his packets, “light-quanta”; the modern term, “photon” was coined by G.N.Lewis . This localized-packet concept explains the photo-electric effect by a quasi-chemical equation:
$`\mathrm{photon}+\mathrm{atom}\mathrm{emitted}\mathrm{\_}\mathrm{electron}+\mathrm{positive}\mathrm{\_}\mathrm{ion}`$ (1)
the energy of the emitted electron being equal to the energy of the absorbed photon minus the energy required to remove the electron from the atom within the surface of the solid forming the photocell: This equation (1) predicts the experimentally observed characteristics of the photo-electric effect:
1. all the emitted electrons have the same kinetic energy when the light is monochromatic,
2. this kinetic energy increases as the frequency of the light increases (as the wavelength decreases),
3. there is a minimum frequency ($`\nu _0`$) of the light (maximum wavelength) below which no electrons are emitted; the energy, $`h\nu _0`$ is the energy required to remove an electron from the surface of the photocell (binding energy of the electron, ionization energy of the atom, within the solid surface),
4. the kinetic energy of the emitted electrons is $`h\nu `$$``$$`h\nu _0`$, which is observed as the voltage of the photo-cell being given by: $`V=(h\nu `$$``$$`h\nu _0)/e`$ where $`e`$ is the electric charge of the electron,<sup>1</sup><sup>1</sup>1$`\nu _0`$ is measured as the reverse (stopping) voltage, $`(h\nu _0)/e`$, needed to just stop the flow of current from the photocell.
5. the electric current generated by the photocell (rate of charge flow $``$ electrons per unit time) is proportional to the radiation intensity absorbed by the cell (equivalent to number of photons per unit area of cell surface absorbed in unit time).
### 1.1 The Photon in Quantum Field Theory and the Bohr Model
Although the quantization of the radiation field in terms of photons of energy $`h\nu `$ became part of the standard language of quantum optics , Einstein’s original concept of the photon as a localized packet of electromagnetic radiation was discarded with the ascendency of quantum mechanics in the mid-1920s.<sup>2</sup><sup>2</sup>2A corollary of Einstein’s localized packet concept of the photon is that much of the cross-sectional area of a beam of light is empty space, the proportion of the beam’s area occupied by packets increasing with increasing light intensity. Theories of the wave function of the photon preclude this localization, and even the quantum theory of the photo-electric effect models the light as a plane wave \[7, pp.215-224\].
In contrast, the Bohr model of the photon predicts the size and shape of photons, and is thus a quantification of Einstein’s localized packet concept. This prediction of size and shape was not an primary objective of the Bohr model of the photon ; the size and shape resulted from imposition of the principle of causality on the chosen solutions of Maxwell’s equations. The result was that a circularly polarized photon is a monochromatic electromagnetic traveling wave confined within a circular ellipsoid of length equal to the wavelength ($`\lambda `$), and of diameter $`\lambda /\pi `$; i.e. an egg-shaped solitary wave propagating along the long axis of the ellipsoid.<sup>3</sup><sup>3</sup>3The ellipsoid is 3 times (accurately $`\pi `$ times) as long as its diameter. This prediction of size and shape was most important because it provided a basis for comparison with experimental observations. The comparison produced agreement between the predicted size and shape and those inferred from several experimental measurements and observations. This agreement makes the Bohr model worthy of serious consideration even though its theoretical basis (a quantized solution of Maxwell’s equations confined by causality) is dissimilar from the widely accepted quantum field theory of light .
This ellipsoidal soliton model of the photon is a Bohr model in the sense that it is a solution of the classical equations of motion that is subsequently quantized. In Bohr’s well-known model of the hydrogen atom the classical equations are Newton’s equations for the motion of an electron within the field of a proton, whereas for the photon (light regarded as electromagnetic radiation) the appropriate classical equations are Maxwell’s equations in vacuum. In Bohr’s model of the hydrogen atom the quantization makes the angular momentum of the electron an integer multiple of Planck’s constant, $`\mathrm{}=h/2\pi `$. In the Bohr model of the photon the quantization of the photon’s angular momentum arises from an appropriately chosen solution of Maxwell’s equations: in addition, the energy of the oscillating electromagnetic field (integrated over the volume of the ellipsoid) is quantized to be $`h\nu `$ \- the known energy of the photon; this quantization fixes the amplitude of the wave, and is analogous to the imposed quantization of angular momentum in Bohr’s model of the hydrogen atom; the analogy extends to quantization of the energy being $`nh\nu `$ with $`n>1`$ representing a multiphoton.<sup>4</sup><sup>4</sup>4Observations indicate that multiphotons have a strong tendency to separate laterally into single photons moving along parallel propagation axes. This instability and the stability of a photon with just one $`h\nu `$ of energy is an outstanding mystery of physics, whose eventually resolution should yield a profound insight into the nature of Planck’s constant.
The full theoretical derivation of the Bohr model is presented in together with supporting experimental evidence; here the theory and experimental support is summarized and augmented by recent ideas pertaining to how a solitary wave can exhibit two-slit interference.
### 1.2 The Bohr Model of the Photon Summarized
The solution of Maxwell’s equations was chosen to be a monochromatic traveling wave having the observed angular momentum of the photon; i.e. a spin of $`\pm \mathrm{}`$; constant parameters multiplying each of these spin states allows for representation of all the known polarization states of light.
The chosen solution of Maxwell’s equations is confined within a finite space-time region by the principle of Special Relativity that causally related events must be separated by time-like intervals. With the idea that a photon is self-causing as it propagates, causality imposes the condition that events within the wave having the same phase must be separated by time-like intervals. In the limit where the interval becomes null (light-like), causality leads to the inference that the length of the photon along its axis of propagation is the wavelength, $`\lambda `$.<sup>5</sup><sup>5</sup>5or equivalently in time, the period of oscillation $`\tau =\nu ^1`$. In addition, for circularly polarized states the causally connected field is contained within a circular ellipsoid with maximum diameter (transverse to the axis of propagation) of $`\lambda /\pi `$; the length of the ellipsoid (along the axis of propagation) is the wavelength.<sup>6</sup><sup>6</sup>6The ellipsoidal soliton can be visualized as an egg, or as an american/rugby football.
This modeling of the photon as an ellipsoidal soliton arises from the imposition of causality upon the solution of Maxwell’s equations (which are linear and homogeneous) whereas non-relativistic solitons arise as solutions of non-linear differential equations .
The size and shape of the soliton allowed for quantization of its energy; the wave’s electromagnetic energy, $`𝐄^2+𝐇^2`$, integrated over the volume of the ellipsoid, was set to $`h\nu `$.<sup>7</sup><sup>7</sup>7Or in general, to $`nh\nu `$, analogous with Bohr’s quantization of angular momentum as $`n\mathrm{}`$. This fixed the amplitude of the wave and led to an expression for the average intensity within the photon-soliton \[8, eqn.57\]:<sup>8</sup><sup>8</sup>8The photon’s intrinsic intensity is not uniform: being proportional to the radius ($`r`$) squared
it is zero on the axis of propagation and maximal at the ellipsoid’s maximum radius of $`r=\lambda /2\pi `$.
$$I_p=\frac{4\pi hc^2}{\lambda ^4}$$
(2)
### 1.3 Experimental Confirmation of the Soliton
Experiment confirms the predicted size, shape and intrinsic intensity of the photon:
* its length of $`\lambda `$ is confirmed by:
+ the generation of laser pulses that are just a few periods long;
+ for the radiation from an atom to be monochromatic (as observed),
the emission must take place within one period, $`\tau `$, ;
+ the sub-picosecond response time of the photoelectric effect ;<sup>9</sup><sup>9</sup>9The predicted absorption time of the ellipsoidal photon is its period of oscillation, $`\tau =1/\nu `$ = the transit time of the ellipsoid past any point in space = the time to enter the surface of the photocell: a few femtoseconds for visible light.
* the diameter of $`\lambda /\pi `$ is confirmed by:
+ the attenuation of direct (undiffracted) transmission of circularly polarized light through slits narrower than $`\lambda /\pi `$: our own measurements of the effective diameter of microwaves \[8, p.166\] confirmed this within the experimental error of 0.5%;
+ the resolving power of a microscope (with monochromatic light) being “a little less than a third of the wavelength”; $`\lambda /\pi `$ is 5% less than $`\lambda /3`$, ;
* The predicted intrinsic intensity (given by eqn.2) is the threshold (minimum) intensity to which a laser beam must be focussed in order to produce multiphoton absoption: two experiments confirming this (one with 650$`nm`$ light , the other with $`\lambda `$=10.5$`\mu m`$) are described in \[8, p.165\].
### 1.4 Solution of Maxwell’s Equations: the Photon’s Wave Function
Maxwell’s equations relate the first derivatives of the six components of the electromagnetic field; they comprise eight partial differential equations which must be satisfied simultaneously.<sup>10</sup><sup>10</sup>10The equations are linear and homogeneous with constant coefficients. The key to finding appropriate solutions, is to differentiate to produce second derivatives followed by elimination of common terms between the resulting equations to yield the result that each Cartesian component of the field ($`E_x,E_y,E_z,H_x,H_y,H_z`$) separately satisfies d’Alembert’s wave equation .<sup>11</sup><sup>11</sup>11This only pertains for the Cartesian components; it does not prevail for the spherical or cylindrical components.
For a wave traveling parallel to the $`z`$-axis at the speed of light, $`c`$, the solution must be any function of $`zct`$ , and if this wave is monochromatic the functional form is simply:<sup>12</sup><sup>12</sup>12$`S(zct)`$ is an eigenfunction of Schrödinger operators: momentum in the direction of propagation, $`\widehat{p}_z=\frac{\mathrm{}}{i}\frac{}{z}`$, with eigenvalue $`h/\lambda `$, and energy, $`\widehat{E}=\frac{\mathrm{}}{i}\frac{}{t}`$, with eigenvalue $`hc/\lambda =h\nu `$, the physically known values for the photon.
$`S(zct)=\mathrm{exp}\{2\pi i(zct)/\lambda \}`$ (3)
When this form is adopted as a factor of the solution, insertion into d’Alembert’s equation causes a complete separation of $`z`$ and $`t`$ from the transverse coordinates ($`x=r\mathrm{cos}\varphi ,y=r\mathrm{sin}\varphi `$),<sup>13</sup><sup>13</sup>13The separation is complete: there is no separation constant between the $`z,t`$ and the $`r,\varphi `$ differential equations. plane polar coordinates ($`r,\varphi `$) being chosen in preference to the Cartesian coordinates ($`x,y`$) in view of the axial symmetry of the direction of propagation.
Separation of the radius, $`r`$, from the polar angle, $`\varphi `$, produces the two ordinary differential equations:
$`{\displaystyle \frac{1}{\mathrm{\Phi }(\varphi )}}{\displaystyle \frac{d^2\mathrm{\Phi }(\varphi )}{d\varphi ^2}}=m^2={\displaystyle \frac{1}{R(r)}}\left\{{\displaystyle \frac{d^2R(r)}{dr^2}}+{\displaystyle \frac{1}{r}}{\displaystyle \frac{dR(r)}{dr}}\right\}`$ (4)
where $`m^2`$ is the real separation constant introduced to separate $`r`$ from $`\varphi `$.
The simplest solution of eqn.(4) is the plane wave ($`m^2=0`$); i.e. $`R(r)`$ and $`\mathrm{\Phi }(\varphi )`$ both being constants.<sup>14</sup><sup>14</sup>14Plane waves are widely used in the quantum field theory of light . However, this solution was rejected as unphysical because light is observed to travel along very narrow beams.<sup>15</sup><sup>15</sup>15A plane wave has field components that have the same value throughout any plane perpendicular to the axis of propagation, and thus it is completely non-localized, contrary to observation that light moves along very narrow beams.
The next simplest solution of eqns.(4) is for $`m^2=1`$: i.e. a factor of $`r`$ or $`1/r`$, with an angular factor of $`\mathrm{exp}\{i(\varphi )\}`$ or $`\mathrm{exp}\{i(\varphi )\}`$. These angular factors are eigenfunctions of the $`z`$-component of angular momentum, $`𝐋_𝐳=\frac{\mathrm{}}{i}\frac{}{\varphi }`$, in Schrödinger quantum mechanics \[17, p.217\], the eigenvalues of $`\pm \mathrm{}`$ being those observed for the spin angular momentum of the photon; thus these solutions for $`m^2`$=$`1`$ are appropriate for the wavefunction of the photon:
$$\psi (r,\varphi ,zct)=(\alpha r+\beta /r)\left(A\mathrm{exp}\{i\varphi \}+B\mathrm{exp}\{i\varphi \}\right)\mathrm{exp}\{2\pi i(zct)/\lambda \}$$
(5)
Having determined this (the form in eqn.5) as the appropriate solution of d’Alembert’s equation, each of the 6 field components ($`E_x,E_y,E_z,H_x,H_y,H_z`$) will have this form, the coefficients ($`\alpha ,\beta ,A,B`$) being different in each component. The relationships between the coefficients of different components were determined by Maxwell’s equations. This produced the inferences:
$`E_z`$ $`=`$ $`H_z=0(\mathrm{no}\mathrm{field}\mathrm{along}\mathrm{the}\mathrm{axis}\mathrm{of}\mathrm{propagation}:\mathrm{a}\mathrm{transverse}\mathrm{wave})`$
$`E_x`$ $`=`$ $`(\alpha r+\beta /r)\left(A\mathrm{exp}\{i\varphi \}+B\mathrm{exp}\{i\varphi \}\right)\mathrm{exp}\{2\pi i(zct)/\lambda \}=\mu _0cH_y`$ (6)
$`E_y`$ $`=`$ $`i(\alpha r\beta /r)\left(A\mathrm{exp}\{i\varphi \}B\mathrm{exp}\{i\varphi \}\right)\mathrm{exp}\{2\pi i(zct)/\lambda \}=\mu _0cH_x`$
Imposition of the causality condition led to the result that if $`A`$ or $`B`$ is zero, then the field must be contained within a circular ellipsoid of length $`\lambda `$ and cross-sectional diameter $`\lambda /\pi `$ \[8, §2.5\].
Since Maxwell’s equations are linear and homogeneous they do not determine the amplitude of the solutions; this was determined by integration of the energy of the wave, $`𝐄^2+𝐇^2`$.<sup>16</sup><sup>16</sup>16This is analogous with Bohr’s quantization of the electron’s angular momentum in his model of the hydrogen atom. This led to the realization that the form $`1/r`$ would cause a divergent contribution to the energy at $`r=0`$, while the form $`r`$ would cause a similar divergence as $`r\mathrm{}`$. Thus, in view of the causality condition limiting the domain of the field to an ellipsoid along the axis of propagation, it was decided to discard the $`1/r`$ form and retain the $`r`$ form in order to produce a finite integrated energy. This discarding of the $`1/r`$ term (i.e. $`\beta `$=0 in eqn.6) was concordant with the need to make the field an eigenfunction of $`L_z`$ \[8, §2.6\].
This normalization of the amplitude of the photon’s field yielded:<sup>17</sup><sup>17</sup>17In the amplitude squared ($`\alpha ^2`$=$`S_0^2`$ in \[8, eqn.47\]) was given as, $`\alpha ^2=64nhc\pi ^4/(ϵ_0\lambda ^6)`$, which corresponds to integration over a cylinder (length $`\lambda `$ and diameter $`\lambda /\pi `$) rather than the ellipsoid; the factor of 120 in eqn.(7) is correct for integration over the ellipsoid; the relation $`A^2+B^2=1`$ was imposed, with $`\alpha ^2`$ determined by the energy integral.
$`A^2+B^2=1\mathrm{and}\alpha ^2=120nhc\pi ^4/(ϵ_0\lambda ^6)`$ (7)
### 1.5 The Soliton’s Evanescent Wave
An evanescent wave outside the ellipsoid is necessary as an adjunct to the theory presented in , because while the relativistic principle of causality confines the wave within the ellipsoid, the radial dependence of the wave within the soliton is simply $`r`$, which is a maximum at the surface of the ellipsoid; physically the wave cannot sharply cut-off to zero at this surface; it must smoothly decay towards zero outside the ellipsoid; an evanescent wave decays in this way \[7, pp.103-108\].
The radial dependence of the evanescent wave is $`1/r`$; i.e. the solution of Maxwell’s equations (eqn.6) with $`\alpha =0`$. The intensity of this wave decreases as $`1/r^2`$ as the distance, $`r`$, from the axis increases.
J.J. Thomson derived the same solution (eqn.6) of Maxwell’s equations in 1924 ; he noted that a radial dependence of $`r`$ is appropriate near $`r=0`$, with $`1/r`$ being appropriate as $`r\mathrm{}`$, but he didn’t pursue his analysis as far as deducing an ellipsoidal soliton, with the wave having the $`r`$ form within the ellipsoid, and the $`1/r`$ form outside the ellipsoid.
The $`r`$ dependence within the ellipsoid and the $`1/r`$ dependence outside the ellipsoid, makes the $`r`$-derivative of the wave discontinuous on the surface of the ellipsoid. While this may appear to be unphysical, it is the same discontinuity exhibited by the gravitational force due to the mass of the Earth: on the assumption of a uniform density, the gravitational force inside the Earth is proportional to the radius, $`r`$, whereas outside the Earth it decreases like $`1/r^2`$ .
### 1.6 Characteristics of the Photon’s Evanescent Wave
The polar components of the evanescent field are given by eqns.(38) of for $`\alpha =0`$ and $`\beta `$ given by eqn.(10), which show that none of these components have any dependence upon the polar angle $`\varphi `$, and that $`E_r`$ and $`H_\varphi `$ are real, while $`H_r`$ and $`E_\varphi `$ are imaginary:
$`E_r={\displaystyle \frac{\beta }{r}}\left[A+B\right]=\mu _0cH_\varphi E_\varphi =i{\displaystyle \frac{\beta }{r}}\left[AB\right]=\mu _0cH_r`$ (8)
Independence of the angle, $`\varphi `$, means that the evanescent wave carries none of the angular momentum of the photon,<sup>18</sup><sup>18</sup>18Because the operator for the $`z`$-component of angular momentum is $`𝐋_𝐳=\frac{\mathrm{}}{i}\frac{}{\varphi }`$. and hence none of its energy; it is a truly evanescent wave \[7, pp.105-108\].
### 1.7 Matching the Soliton and Evanescent Waves
While the gradient of the wave has a cusp at $`r=\lambda /(2\pi )`$, the amplitude must be continuous at $`r=\lambda /(2\pi )`$; equating of the soliton and evanescent wave amplitudes at $`r=\lambda /(2\pi )`$ produces:
$`\alpha r=\beta /r\mathrm{for}\mathrm{r}=\lambda /(2\pi )`$ (9)
and since $`\alpha ^2`$ is given by eqn.(7) it follows that:
$`\beta ^2=[\lambda /(2\pi )]^4\times 120nhc\pi ^4/(ϵ_0\lambda ^6)=7.5nhc/(ϵ_0\lambda ^2)`$ (10)
#### Orthogonality of the Radial Gradients
The radial gradient of the soliton wave is simply the normalization constant, $`\alpha `$, while that of the evanescent wave is $`\beta /r^2`$. Thus at the cusp where the two waves join (at $`r=\lambda /(2\pi )`$) the ratio of these gradients is:
$`\mathrm{ratio}\mathrm{of}\mathrm{gradients}={\displaystyle \frac{\beta }{\alpha r^2}}=1\mathrm{at}r=\lambda /(2\pi )`$ (11)
Thus where the soliton and evanescent waves meet (at $`r=\lambda /(2\pi )`$) they are orthogonal to each other - independent of the wavelength, $`\lambda `$.
The above matching of the soliton and evanescent waves was made at the soliton’s maximum diameter of $`\lambda /\pi `$; this raises the question of their matching at values of $`z`$ other than $`z=0`$; i.e. at other points on the ellipse:
$`\left(2\pi r\right)^2+\left(2z\right)^2=\lambda ^2`$
$`\mathrm{i}.\mathrm{e}.\mathrm{when}r`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}\sqrt{(\lambda )^2(2z)^2}\mathrm{for}{\displaystyle \frac{\lambda }{2}}<z<+{\displaystyle \frac{\lambda }{2}}`$ (12)
It might appear natural to apply the matching condition of eqn.(9) for all values of $`r`$ specified in eqn.(12) to produce:
$`\beta ^2`$ $`=`$ $`\left[{\displaystyle \frac{1}{2\pi }}\sqrt{(\lambda )^2(2z)^2}\right]^4\times 120nhc\pi ^4/(ϵ_0\lambda ^6)`$ (13)
$`=`$ $`\left[(\lambda )^2(2z)^2\right]^2\times 7.5nhc/(ϵ_0\lambda ^6)`$
This would have the effect of making the amplitude of the evanescent wave, $`\beta `$, smaller as $`z`$ changes from $`z=0`$ to $`z=\pm \frac{\lambda }{2}`$, with $`\beta `$ being zero at these limits (the ends of the ellipsoid).
However, this conjecture would make $`\beta `$ a function of $`z`$ (as in eqn.13) rather than a constant, and hence the evanescent field (eqn.6 for $`\alpha =0,\beta 0`$) would no longer be a a solution of Maxwell’s equations. The resolution of this physical vs. mathematical paradox may be found within the framework of General Relativity, in which the photon’s local energy produces a non-Lorentzian metric.
### 1.8 Diffraction and Interference
The evanescent wave is believed to be responsible for the phenomena of diffraction and interference. As a photon-soliton passes close to the edge of, or through a slit in, a material obstacle placed within the beam of light, the interaction between the electrons within the obstacle and the photon’s evanescent wave will cause its path to bend as it passes by, the angle of bending (diffraction) being dependent upon the impact parameter of the soliton’s axis with the edge or slit.
Double slit interference can be understood by the soliton itself (like the $`C_{60}`$ molecules in Zeilinger’s experiment ) going through one slit or the other, while its evanescent wave extends over both slits. The evanescent wave is like a classical continuous wave in extending throughout all space, and hence the interference minima and maxima will appear at the same positions as predicted by Huygen’s theory. However, the soliton model predicts that:
* the individual photons will arrive at local positions in the detection plane, whereas the classical continuous wave model predicts a uniformly visible interference pattern: that the former (rather than the latter) is actually observed supports the soliton model ;
* the visibility of the interference pattern<sup>19</sup><sup>19</sup>19Visibility, $`V`$, is defined by: $`V=\frac{I_{\mathrm{max}}I_{\mathrm{min}}}{I_{\mathrm{max}}+I_{\mathrm{min}}}`$, $`I_{\mathrm{max}}`$ and $`I_{\mathrm{min}}`$ being the measured intensities at the interference maxima and minima respectively; it has the range: $`0V1`$. will decrease with slit separation (because the intensity of the evanescent wave decreases like $`1/r^2`$, $`r`$ being the distance from the soliton’s axis of propagation), whereas the classical continuous wave model predicts a visibility independent of slit separation. This seems not to have been investigated experimentally .
A double-slit experiment by Alkon exhibits the expected interference pattern even though the individual photons are constrained to pass through one slit or the other by an opaque barrier extending from the source (a laser) up to the mid-point between the slits.<sup>20</sup><sup>20</sup>20Alkon’s experiment is the experimental proof that the continuous wave concept that “the photon goes through both slits and interferes with itself” is not correct. This experiment demonstrates that the particle-like photon (the Bohr model soliton) passes through one slit or the other, and yet its passage through this slit (and the subsequent diffraction) is affected by the presence of the other slit; this effect of the other open slit is evidence for the existence of the evanescent wave surrounding the soliton.<sup>21</sup><sup>21</sup>21Interaction between the evanescent waves of collaterally moving photon-solitons could be the cause of the very small (but finite) divergence of a laser beam \[23, p.6\].
A causal model of diffraction has been proposed by Gryzinski ; it is based upon the photon being a particle-like (localized) electromagnetic wave that interacts with the array of positive atomic nuclei and negative electrons within a solid, as it passes:
* through a crystal (Bragg diffraction of X-rays), or
* adjacent to an edge of a sheet of the solid (an edge of a slit).
Gryzinsky’s model of diffraction does not specify the size or shape of the soliton, but it quantitatively explains both Bragg diffraction and double-slit interference; his concept of the latter is that while the localized photon goes through one slit, its wave extends to the other slit. His theory is concordant with the Bohr model’s evanescent wave, specifically because his localized model involves the concept that “the photon’s electric field decreases when distance \[from its center\] increases”.
Gryzinsky pertinently cites Zeilinger’s observation that each photon manifests its particle (localized) nature in each detection event: the distribution of detection events<sup>22</sup><sup>22</sup>22attributed in the continuous wave model to the wave going through both slits and self-interfering only becomes manifest after a large number ($`10^4`$) of detection events have been recorded ; each photon detection is a localized event.
The evanescent wave explanation for diffraction and interference is not readily invoked for the Mach-Zender type of interferometer, because the two alternative paths for the photon are typically separated by distances over which the evanescent wave’s intensity would have become negligible; a small difference (of the order of the wavelength) between the lengths of the two paths determines the observed interference pattern. This observation requires further theoretical explanation. |
warning/0506/math0506567.html | ar5iv | text | # A geometric classification of immersions of 3-manifolds into 5-space
## 1. Introduction
Hirsch-Smale theory reduces the problem of regular homotopy classification of immersions to homotopy theory. However, this homotopy theoretic problem is usually hard to deal with. In the case of immersions of oriented 3-manifolds into $`^5`$ this homotopy theoretic problem was solved by Wu using algebraic topological methods (also see ).
However, it remains a problem to determine the regular homotopy class of a given immersion from its geometry. By geometry we mean the structure of the ”singularities” of the map. For example, double points are such ”singularities”, and indeed, Smale showed that for $`n>1`$ the regular homotopy class of an immersion $`S^n^{2n}`$ is completely determined by the number of its double points (modulo 2 if $`n`$ is odd). A similar classification was carried out by Ekholm for immersions of $`S^k`$ into $`^{2k1}`$ for $`k4`$.
The whole picture changes when we consider immersions of $`S^3`$ into $`^5`$. Hughes and Melvin showed that there are infinitely many embeddings $`S^3^5`$ that are pairwise not regularly homotopic to each other. Therefore one can not determine the regular homotopy class from the ”singularities” since an embedding has no such. Ekholm and Szűcs came over this problem using ”singular Seifert surfaces” bounded by the immersions. For an immersion $`f:M^3^5`$ a singular Seifert surface is a generic map $`F:W^4^5`$ of a compact orientable manifold $`W^4`$ with boundary $`M^3`$ such that $`F=f`$. In it is shown that for $`M^3=S^3`$ the Smale invariant of $`f`$ can be computed from the singularities of $`F`$. Later Saeki, Szűcs and Takase generalized these results for immersions $`f:M^3^5`$ with trivial normal bundle (for oriented $`M^3`$). The invariant introduced in corresponds to the 3-dimensional obstruction to a regular homotopy between two such immersions. Our present paper generalizes the results of to arbitrary immersions $`f:M^3^5`$.
We will consider the set $`\text{Imm}(M^3,^5)_\chi `$ of immersions with fixed normal Euler class $`e(\nu _f)=\chi H^2(M^3;)`$ and construct a $`_{2d(\chi )}`$-valued regular homotopy invariant $`i`$ for this set of immersions, where $`d(\chi )`$ denotes the divisibility of $`\chi `$. The construction of the invariant $`i`$ will also make use of a singular Seifert surface $`F`$. In $`F`$ had to be an immersion near the boundary, but we (have to and) will allow arbitrary generic maps. If $`\chi =0`$ and $`F`$ is an immersion near the boundary then the construction of the invariant $`i`$ agrees with the one introduced in . We will also show that whenever $`f,g:M^3^5`$ are regularly homotopic on a neighborhood of the 2-skeleton of $`M^3`$ then $`i(f)=i(g)`$ iff $`f`$ and $`g`$ are regularly homotopic. This shows that $`i`$ corresponds to the 3-dimensional obstruction to a regular homotopy between $`f`$ and $`g`$. (Note that there is an invariant which determines the regular homotopy class of the restriction of an immersion to a neighborhood of the 2-skeleton of $`M^3`$. This invariant was called the Wu invariant in , see below.)
Regular homotopy classes of immersions of oriented 3-manifolds into $`^5`$ endowed with the connected sum operation form a semigroup whose structure we will also determine. Finally, an exact sequence will be defined that relates $`\text{Imm}[M^3,^5]`$ to $`\text{Imm}[M^3,^6]`$ and $`[M^3,S^2]`$.
## 2. Preliminaries
First let us recall the result of Wu that classifies immersions of an oriented 3-manifold $`M^3`$ into $`^5`$ up to regular homotopy.
###### Theorem 2.1.
The normal Euler class $`\chi `$ of an immersion $`f:M^3^5`$ is of the form $`2c`$ for some $`cH^2(M^3;)`$ and for any $`cH^2(M^3;)`$ there is an immersion $`f`$ such that $`\chi =2c`$. Furthermore,
$$\text{Imm}[M^3,^5]_\chi \underset{cH^2(M^3;),\mathrm{\hspace{0.17em}2}c=\chi }{}H^3(M^3;)/(2\chi H^1(M^3;)),$$
where $`\text{Imm}[M^3,^5]_\chi `$ is the set of regular homotopy classes of immersions with normal Euler class $`\chi H^2(M^3;)`$ and $``$ represents the cup product, moreover the symbol $``$ denotes a bijection.
###### Remark 2.2.
For $`\chi H^2(M^3;)`$ let $`d(\chi )`$ denote the divisibility of $`\chi `$, so that $`\chi `$ equals $`d(\chi )`$ times a primitive class in $`H^2(M^3;)`$ modulo torsion, and $`d(\chi )=0`$ if $`\chi `$ is of finite order. Then Poincaré duality implies that
$$H^3(M^3;)/(2\chi H^1(M^3;))_{2d(\chi )}.$$
If $`f`$ is an immersion of $`M^3`$ into $`^5`$ with normal Euler class $`\chi `$ then let us introduce the notation $`d(f)`$ for $`d(\chi )`$.
###### Notation 2.3.
For $`\chi H^2(M^3;)`$ let $`\mathrm{\Gamma }_2(\chi )`$ denote the set $`\{cH^2(M^3;):2c=\chi \}`$. Throughout this paper we will use the notation $`M_{}^3`$ for the punctured 3-manifold $`M^3D^3`$, where $`D^3M^3`$ is a closed 3-disc. Then the 2-skeleton $`\text{sk}_2(M^3)`$ is a deformation retract of $`M_{}^3`$.
Theorem 2.1 can also be applied to the open manifold $`M_{}^3`$. Since $`H^3(M_{}^3;)=0`$ we obtain a bijection
$$\overline{c}:\text{Imm}[M_{}^3,^5]_\chi \mathrm{\Gamma }_2(\chi ).$$
Thus for an immersion $`f:M^3^5`$ the invariant $`c(f)=\overline{c}(f|M_{}^3)\mathrm{\Gamma }_2(\chi )`$ describes the regular homotopy class of $`f|M_{}^3`$. Following we will call $`c(f)`$ the Wu invariant of the immersion $`f`$.
To get a complete description of $`\text{Imm}[M^3,^5]_\chi `$ we will construct a $`_{2d(\chi )}`$-valued invariant $`i`$ such that the map
$$(c,i):\text{Imm}[M^3,^5]_\chi \mathrm{\Gamma }_2(\chi )\times _{2d(\chi )}$$
will be a bijection. The invariant $`i`$ is constructed in a geometric manner and is an extension of the invariant defined in for $`\chi =0`$.
Next let us recall Theorem 1.1(a) in . Let $`f:S^3^5`$ be an immersion and $`V^4`$ an arbitrary compact oriented 4-manifold with $`V^4=S^3`$. The map $`f`$ extends to a generic map $`F:V^4^5`$ which has no singular points near the boundary $`V^4`$ since the normal bundle $`\nu _f`$ of $`f`$ is trivial. This map $`F`$ has isolated cusps, each one having a sign. Let us denote by $`\mathrm{\#}\mathrm{\Sigma }^{1,1}(F)`$ their algebraic number and let $`\mathrm{\Omega }(f)`$ be the Smale invariant of $`f`$. The following formula was proved in .
###### Theorem 2.4.
$$\mathrm{\Omega }(f)=\frac{1}{2}(3\sigma (V^4)+\mathrm{\#}\mathrm{\Sigma }^{1,1}(F)).$$
The proof of this theorem relies on the following proposition (, Lemma 3).
###### Lemma 2.5.
Let $`X^4`$ be a closed oriented 4-manifold and $`g:X^4^5`$ a generic map. Then $`3\sigma (X^4)+\mathrm{\#}\mathrm{\Sigma }^{1,1}(g)=0`$.
For the sake of completeness we will also recall from the definition of the invariant $`i`$ for immersions with trivial normal bundle. First we need a preliminary definition.
###### Definition 2.6.
Let $`M^3`$ be a closed oriented 3-manifold. We denote by $`\alpha (M^3)`$ the dimension of the $`_2`$ vector space $`\tau H_1(M^3;)_2`$, where $`\tau H_1(M^3;)`$ is the torsion subgroup of $`H_1(M^3;)`$.
###### Definition 2.7.
Let $`f:M^3^5`$ be an immersion with trivial normal bundle. Let $`W^4`$ be any compact oriented 4-manifold with $`W^4=M^3`$ and $`F:W^4^5`$ a generic map nonsingular near the boundary such that $`F|W^4=f`$. (We can choose such a generic map $`F`$ since $`f`$ is an immersion with trivial normal bundle.) Denote the algebraic number of cusps of $`F`$ by $`\mathrm{\#}\mathrm{\Sigma }^{1,1}(F)`$. Then let
$$i(f)=\frac{3}{2}(\sigma (W^4)\alpha (M^3))+\frac{1}{2}\mathrm{\#}\mathrm{\Sigma }^{1,1}(F).$$
It is proved in that $`i(f)`$ is always an integer and a regular homotopy invariant.
In the following sections we will extend the above regular homotopy invariant $`i`$ to arbitrary immersions. If $`f:M^3^5`$ has non-trivial normal bundle then we have to give up the assumption that the singular Seifert-surface $`F`$ is an immersion near the boundary. Thus we will use an arbitrary generic map $`F:W^4^5`$ such that $`F=f`$. The singular set $`\mathrm{\Sigma }^1(F)`$ of such an $`F`$ is a 2-dimensional submanifold of $`W^4`$ with boundary $`C(F)=\mathrm{\Sigma }^1(F)M^3`$. If we orient $`\mathrm{ker}(dF)|C(F)`$ so that it points into $`W^4`$ and project it into $`TM^3`$ then we obtain a normal field $`\nu (F)`$ along $`C(F)`$. We will define the rotation of $`\nu (F)`$ around $`C(F)`$ modulo $`4d(f)`$ and denote this by $`R(F)`$. The double of the extended invariant will be defined to be
$$I(f)=3(\sigma (W^4)\alpha (M^3))+\mathrm{\#}\mathrm{\Sigma }^{1,1}(F)+R(F)_{4d(f)}.$$
We will show that $`I(f)`$ is always even, thus it defines an element $`i(f)_{2d(f)}`$ using the natural embedding $`_{2d(f)}_{4d(f)}`$.
## 3. Rotation
Throughout this paper $`M^3`$ will denote a fixed closed connected and oriented 3-manifold.
###### Notation 3.1.
A pair $`(C,\nu )`$ will always stand for an oriented 1-dimensional submanifold $`C`$ of $`M^3`$ and a nowhere vanishing normal field $`\nu `$ along $`C`$.
###### Definition 3.2.
Let $`\chi H^2(M^3;)`$ and let $`C_0`$ and $`C_1`$ be 1-dimensional oriented submanifolds of $`M^3`$ with normal fields $`\nu _0`$ and $`\nu _1`$ such that $`PD[C_0]=PD[C_1]=\chi `$. (Here $`PD`$ denotes Poincaré duality.) Then we can define the *rotation difference* $`\text{rd}((C_0,\nu _0),(C_1,\nu _1))_{2d(\chi )}`$ of $`(C_0,\nu _0)`$ and $`(C_1,\nu _1)`$ as follows. Since $`[C_0]=[C_1]`$ and
$$H_1(M^3;)H^2(M^3;)[M^3,P^{\mathrm{}}],$$
there exists an oriented cobordism $`K^2M^3\times I`$ between $`C_0M^3\times \{0\}`$ and $`C_1M^3\times \{1\}`$. Let $`\nu `$ be a generic normal field along $`K^2`$ that extends $`\nu _0`$ and $`\nu _1`$. Then a sign can be given to each zero of $`\nu `$ since $`M^3`$ is oriented. Now we define $`\text{rd}((C_0,\nu _0),(C_1,\nu _1))`$ to be the algebraic number of zeroes of $`\nu `$ modulo $`2d(\chi )`$. Equivalently, $`\text{rd}((C_0,\nu _0),(C_1,\nu _1))`$ is the self intersection of $`K`$ in $`M^3\times I`$ modulo $`2d(\chi )`$ if perturbed in the direction of $`\nu `$.
###### Remark 3.3.
The rotation difference is the obstruction to the existence of a framed cobordism between the framed submanifolds $`(C_0,\nu _0)`$ and $`(C_1,\nu _1)`$ of $`M^3`$. Using the Pontrjagin construction this corresponds to the obstruction to a homotopy between two maps of $`M^3`$ to $`S^2`$. This situation was first examined in . It is easy to see that $`\text{rd}((C_0,\nu _0),(C_1,\nu _1))=0`$ iff $`(C_0,\nu _0)`$ and $`(C_1,\nu _1)`$ are framed cobordant. Thus we obtain a bijection
$$[M^3,S^2]\underset{\chi H^2(M^3;)}{}H^3(M^3;)/2\chi H^1(M^3;).$$
###### Proposition 3.4.
In Definition 3.2 above the rotation difference is well defined, i.e., it does not depend on the choice of $`K`$ and $`\nu `$.
###### Proof.
Let $`K`$, $`\nu `$ and $`K^{}`$, $`\nu ^{}`$ be as in Definition 3.2. We glue together $`M^3\times I`$ and $`M^3\times I`$ along their boundaries so that we obtain the double $`D(M^3\times I)=M^3\times S^1`$. Place $`K`$ into the half of $`M^3\times S^1`$ corresponding to $`M^3\times I`$ and $`K^{}`$ into the other half. Then we obtain a closed oriented surface $`F=KK^{}`$ in $`M^3\times S^1`$ and a normal field $`\mu =\nu \nu ^{}`$ along $`F`$. Since $`H^{}(S^1;)`$ is a torsion free $``$-module we can apply Künneth’s theorem and we get that
$$H^2(M^3\times S^1;)H^1(M^3;)H^1(S^1;)H^2(M^3;)H^0(S^1;).$$
Thus the Poincaré dual of $`F`$ can be written in the form
$$PD[F]=x\times \alpha +y\times 1H^2(M^3\times S^1;),$$
where $`xH^1(M^3;)`$ and $`yH^2(M^3;)`$, moreover $`\alpha `$ denotes the generator of $`H^1(S^1;)`$ and $`1`$ the generator of $`H^0(S^1;)`$ given by the orientation of $`S^1`$. Note that for $`1S^1`$ the dual class of $`M^3\times \{1\}M^3\times S^1`$ is
$$PD[M^3\times \{1\}]=PD[M^3]\times PD[\{1\}]=1\times \alpha H^1(M^3\times S^1;).$$
Moreover,
$$PD[F(M^3\times \{1\})]=PD[C_1\times \{1\}]=PD[C_1]\times PD[\{1\}]=\chi \times \alpha .$$
On the other hand
$$PD[F(M^3\times \{1\})]=PD[F]PD[M^3\times \{1\}]=(x\times \alpha +y\times 1)(1\times \alpha )=x\times \alpha ^2+y\times \alpha .$$
Since $`\alpha ^2=0`$ we get that $`y\times \alpha =\chi \times \alpha `$. Using Künneth’s theorem again we obtain the equality $`y=\chi `$. Thus we get that
$$PD[F]PD[F]=(x\times \alpha +\chi \times 1)^2=(2x\chi )\times \alpha $$
since $`\alpha ^2=\chi ^2=0`$ and $`x\chi =\chi x`$ because the degree of $`\chi `$ is 2. So the self intersection of $`F`$ in $`M^3\times S^1`$ equals $`(2\chi x)\times \alpha ,[M^3\times S^1]=2\chi x,[M^3]2d(\chi )`$. If we perturb $`F`$ in the direction of $`\mu `$ we get that the self intersection of $`K`$ with respect to $`\nu `$ equals the self intersection of $`K^{}`$ with respect to $`\nu ^{}`$ modulo $`2d(\chi )`$. ∎
###### Proposition 3.5.
If $`[C_0]=[C_1]=[C_2]H_1(M^3,)`$ then
$$\text{rd}((C_0,\nu _0),(C_1,\nu _1))+\text{rd}((C_1,\nu _1),(C_2,\nu _2))=\text{rd}((C_0,\nu _0),(C_2,\nu _2)).$$
###### Definition 3.6.
For each $`aH_1(M^3;)`$ fix a pair $`(C_a,\nu _a)`$ such that $`[C_a]=a`$. Then for $`[C]=a`$ let $`r(C,\nu )=\text{rd}((C,\nu ),(C_a,\nu _a))`$.
###### Corollary 3.7.
If $`[C_0]=[C_1]`$ then $`r(C_0,\nu _0)r(C_1,\nu _1)=\text{rd}((C_0,\nu _0),(C_1,\nu _1))`$.
###### Definition 3.8.
We can define the mod 2 rotation difference $`\text{rd}_2((C_0,\nu _0),(C_1,\nu _1))`$ for unoriented $`C_0`$ and $`C_1`$ just as in Definition 3.2 but allowing the cobordism $`K`$ to be non-orientable and counting the self intersection of $`K`$ in $`M\times I`$ only modulo 2. The proof that this is well defined is analogous to the oriented case. It is clear that the epimorphism $`_{2d(\chi )}_2`$ takes rd to $`\text{rd}_2`$. The mod 2 rotation $`r_2`$ is defined just like $`r`$.
Unfortunately we will have to lift the invariants rd and $`r`$ to $`_{4d(\chi )}`$. To be able to do this we need more structure on $`M^3`$ then just a framed submanifold. We will use this additional structure to restrict the homology class of the cobordism $`K`$ so that the surface $`F`$ in the proof of Proposition 3.4 will represent an even homology class and thus $`x`$ will always be even (since $`\chi `$ is even). So the self intersection of $`F`$ will be divisible by $`4d(\chi )`$ instead of just $`2d(\chi )`$.
###### Notation 3.9.
Fix a cohomology class $`\chi H^2(M^3;)`$. Let $`\epsilon _M^3`$ denote the 3-dimensional trivial bundle over $`M^3`$ and let $`t,v\mathrm{\Gamma }(\epsilon _M^3)`$ be two generic non-zero sections of $`\epsilon _M^3`$. Furthermore, suppose that the 2-dimensional oriented subbundle $`t^{}<\epsilon _M^3`$ has Euler class $`\chi `$. If we project $`v`$ into $`t^{}`$ we obtain a section $`w\mathrm{\Gamma }(t^{})`$ that vanishes along a curve $`CM^3`$ and we orient $`C`$ so that $`PD[C]=e(t^{})`$. In particular, $`t`$ and $`v`$ are linearly dependent exactly at the points of $`C`$. Finally let $`\nu `$ be a non-zero normal field along $`C`$. In the future we will denote such a structure on $`M^3`$ by a quadruple $`(C,\nu ,t,v)`$ and the set of these structures by $`N(M^3,\chi )`$.
###### Remark 3.10.
Since $`PD[C]|_2=w_2(\epsilon _M^3)=0H^2(M^3;_2)`$, the cohomology class $`\chi =PD[C]`$ is of the form $`2c`$ for some $`cH^2(M^3;)`$. This can be seen from the long exact sequence associated to the coefficient sequence $`_2`$. Thus $`N(M^3,\chi )=\mathrm{}`$ if $`\chi `$ is not of the form $`2c`$.
###### Definition 3.11.
Suppose that $`a_0=(C_0,\nu _0,t_0,v_0)`$ and $`a_1=(C_1,\nu _1,t_1,v_1)`$ are elements of $`N(M^3,\chi )`$, where $`\chi =2c`$. Then we will define their rotation difference $`\text{Rd}(a_0,a_1)_{4d(\chi )}`$ as follows. We will consider $`a_i`$ to be in $`N(M^3\times \{i\},\chi )`$ for $`i=0,1`$. Let $`t,v\mathrm{\Gamma }(\epsilon _{M\times I}^3)`$ be generic non-zero sections extending $`t_i`$ and $`v_i`$ for $`i=0,1`$. Denote by $`K`$ the 2-dimensional submanifold of $`M^3\times I`$ where $`t`$ and $`v`$ are linearly dependent. Let $`w`$ denote the projection of $`v`$ into the 2-dimensional oriented subbundle $`t^{}<\epsilon _{M\times I}^3`$. Then $`w`$ is zero exactly at the points of $`K`$, thus it defines an orientation of $`K`$. With this orientation $`K`$ is an oriented cobordism between $`C_0`$ and $`C_1`$. Let $`\nu `$ denote a normal field of $`K`$ that extends both $`\nu _0`$ and $`\nu _1`$. Now we define $`\text{Rd}(a_0,a_1)`$ to be the algebraic number of zeroes of $`\nu `$ modulo $`4d(\chi )`$. Equivalently, $`\text{Rd}(a_0,a_1)`$ is the self intersection of $`K`$ in $`M^3\times I`$ modulo $`4d(\chi )`$ if perturbed in the direction of $`\nu `$.
###### Proposition 3.12.
In Definition 3.11 the rotation difference is well defined. I.e., it does not depend on the extensions $`t,v`$ and $`\nu `$.
###### Proof.
Let $`t,v,\nu `$ and $`t^{},v^{},\nu ^{}`$ be as in Definition 3.11. The sections $`t,v`$ are linearly dependent over $`K`$ and $`t^{},v^{}`$ are dependent over $`K^{}`$. Just as in the proof of Proposition 3.4 we will place $`K,\nu `$ and $`K^{},\nu ^{}`$ in the two halves of the double $`D(M^3\times I)=M^3\times S^1`$ and place $`t,v`$ and $`t^{},v^{}`$ in the two halves of the trivial bundle $`\epsilon _{M\times S^1}^3`$. Let $`F`$ denote the oriented surface $`KK^{}`$ and by $`\mu `$ the normal field along $`F`$ obtained from $`\nu `$ and $`\nu ^{}`$. Moreover, let $`T=tt^{}`$ and $`V=vv^{}`$. Then $`T,V\mathrm{\Gamma }(\epsilon _{M\times S^1}^3)`$ are linearly dependent exactly over $`F`$, thus $`PD[F]|_2=w_2(\epsilon _{M\times S^1}^3)=0H^2(M^3\times S^1;_2)`$. Using the coefficient sequence $`_2`$ we get that $`PD[F]`$ is of the for $`2b`$ for some $`bH^2(M^3\times S^1;)`$. Since $`PD[F]`$ is of the form $`x\times \alpha +\chi \times 1`$ where $`\chi =2c`$ we get that there exists an element $`zH^1(M^3;)`$ such that $`x=2z`$. Thus $`PD[F]^2=(4z\chi )\times \alpha `$ which implies that the self intersection of $`F`$ is divisible by $`4d(\chi )`$. If we perturb $`F`$ in the direction of $`\mu `$ we get that the self intersection of $`K`$ with respect to $`\nu `$ equals the self intersection of $`K^{}`$ with respect to $`\nu ^{}`$ modulo $`4d(\chi )`$. ∎
###### Remark 3.13.
The surface $`K`$ represents the dual of the Stiefel-Whitney class of the bundle $`\epsilon _{M\times I}^3`$ relative to the sections $`t_i,v_i`$ given over $`M^3\times \{0,1\}`$. I.e.,
$$PD[K]|_2=w_2(\epsilon _{M\times I}^3;t_i,v_i)H^2(M^3\times I,((M^3C_0)\times \{0\})((M^3C_1)\times \{1\});_2)$$
since $`v`$ and $`t`$ are linearly independent over $`((M^3C_0)\times \{0\})((M^3C_1)\times \{1\})`$. Using Lefschetz duality we get that the relative homology class $`[K]|_2H^2(M^3\times I,C_0\times \{0\}C_1\times \{1\};_2)`$ is independent of the choice of $`t`$ and $`v`$. If we choose a simplicial subdivision of $`M^3`$ so that $`\text{sk}_1(M^3)C_i=\mathrm{}`$ for $`i=0,1`$ then $`w_2`$ is the obstruction to extending the map $`(t_i,v_i):\text{sk}_1(M^3\times \{0,1\})V_2(^3)`$ to $`\text{sk}_2(M^3\times I)`$. So the homology class $`[K]|_2`$ and thus $`\text{Rd}(a_0,a_1)`$ depends only on the homotopy class of the map
$$(t_i,v_i)|\text{sk}_1(M^3):\text{sk}_1(M^3)V_2(^3)$$
for $`i=0,1`$. For the sake of completeness we note that if the extension $`t`$ is given then
$$PD[K]=e(t^{};w_i)H^2(M^3\times I,(M^3C_0)\times \{0\}(M^3C_1)\times \{1\};).$$
So we have obtained the following proposition.
###### Proposition 3.14.
Suppose that $`(C_0,\nu _0)`$ and $`(C_1,\nu _1)`$ are framed submanifolds of $`M^3`$ and let $`a_0,b_0,a_1,b_1N(M^3,\chi )`$ be of the form $`a_i=(C_i,\nu _i,t_i^a,v_i^a)`$ and $`b_i=(C_i,\nu _i,t_i^b,v_i^b)`$ for $`i=0,1`$. Moreover, suppose that $`\text{sk}_1(M^3)C_i=\mathrm{}`$ and $`(t_i^a,v_i^a)|\text{sk}_1(M^3)`$ is homotopic to $`(t_i^b,v_i^b)|\text{sk}_1(M^3)`$ as maps into $`V_2(^3)`$ for $`i=0,1`$. Then the following equality holds:
$$\text{Rd}(a_0,a_1)=\text{Rd}(b_0,b_1).$$
###### Proof.
Let $`t^a`$ and $`v^a`$ be generic extensions of $`t_i^a`$, respectively $`v_i^a`$ over $`M^3\times I`$ and denote by $`K^a`$ the submanifold of $`M^3\times I`$ where $`t^a`$ and $`v^a`$ are linearly dependent. We obtain the sections $`t^b`$ and $`v^b`$ of $`\epsilon _{M\times I}^3`$ and the submanifold $`K^bM^3\times I`$ in a similar way. Then, according to Remark 3.13, we get that
$$PD[K^a]|_2=w_2(\epsilon _{M\times I}^3;t_i^a,v_i^a)=w_2(\epsilon _{M\times I}^3;t_i^b,v_i^b)=PD[K^b]|_2,$$
since $`w_2`$ is the obstruction to extending a map into $`V_2(^3)`$ from $`\text{sk}_1(M^3\times I)`$ to $`\text{sk}_2(M^3\times I)`$ and $`(t_i^a,v_i^a)|\text{sk}_1(M^3\times \{i\})`$ is homotopic to $`(t_i^b,v_i^b)|\text{sk}_1(M^3\times \{i\})`$. Thus if $`F`$ denotes the submanifold of $`M^3\times S^1=D(M^3\times I)`$ obtained by piecing together $`K^a`$ and $`K^b`$ we get that $`PD[F]|_2=w_2(\epsilon _{M\times S^1}^3)=0`$, so we can proceed as in the proof of Proposition 3.12. ∎
###### Proposition 3.15.
If $`a_i=(C_i,\nu _i,t_i,v_i)N(M^3,\chi )`$ for $`i=0,1`$ then
$$\text{rd}((C_0,\nu _0),(C_1,\nu _1))\text{Rd}(a_0,a_1)mod2d(\chi ).$$
###### Proposition 3.16.
If $`a_0,a_1,a_2N(M^3,\chi )`$ then
$$\text{Rd}(a_0,a_1)+\text{Rd}(a_1,a_2)=\text{Rd}(a_0,a_2).$$
###### Definition 3.17.
For each $`\chi H^2(M^3;)`$ of the form $`\chi =2c`$ fix an element $`a_\chi N(M^3,\chi )`$. Then for each $`aN(M^3,\chi )`$ define the dotation $`R(a)_{4d(\chi )}`$ to be $`\text{Rd}(a,a_\chi )`$.
###### Corollary 3.18.
If $`a_0,a_1N(M^3,\chi )`$ then $`\text{Rd}(a_0,a_1)=R(a_0)R(a_1)`$.
## 4. The orientation of $`\mathrm{\Sigma }^1`$
Now let us recall a special case of Lemma 6.1 of . Let $`F:W^4^5`$ be a generic map of a compact orientable manifold. Then the singularity set $`\mathrm{\Sigma }(F)`$ of $`F`$ is a 2-dimensional submanifold of $`W^4`$ which is not necessarily orientable.
###### Lemma 4.1.
The line bundles $`det(T\mathrm{\Sigma }(F))`$ and $`\mathrm{ker}(dF)`$ over $`\mathrm{\Sigma }(F)`$ are isomorphic.
###### Definition 4.2.
Let $`\pi `$ denote the projection of $`^{m+1}`$ onto $`^m`$. A map $`f:N^n^m`$ is called *prim* if there exists an immersion $`f^{}:N^n^{m+1}`$ such that $`\pi f^{}=f`$.
###### Corollary 4.3.
If $`F:W^4^5`$ is a generic prim map then $`\mathrm{\Sigma }(F)W^4`$ is an orientable surface.
###### Proof.
Let $`s`$ denote the sixth coordinate function of $`F^{}`$, i.e., $`F^{}=(F,s)`$. Since $`F^{}`$ is non-singular, the function $`s`$ is non-degenerate along $`\mathrm{ker}(dF)`$. Thus we can orient $`\mathrm{ker}(dF)`$ so that the derivative of $`s`$ in the positive direction of $`\mathrm{ker}(dF)`$ is positive. But the orientability of $`\mathrm{ker}(dF)`$ implies the orientability of $`\mathrm{\Sigma }(F)`$ by Lemma 4.1. ∎
The following definition, motivated by Corollary 4.3, gives an explicit isomorphism $`\mathrm{\Psi }`$ between $`\mathrm{ker}(dF)`$ and $`det(T\mathrm{\Sigma }(F))`$.
###### Definition 4.4.
Let $`W^4`$ be a compact oriented manifold with possibly non-empty boundary and let $`F:W^4^5`$ be a generic map. For $`p\mathrm{\Sigma }(F)`$ choose a small neighborhood $`U_pW^4`$ of $`p`$ in which $`\mathrm{ker}(dF)`$ is orientable. Put $`F_p=F|U_p`$ and choose an orientation $`o_p`$ of $`\mathrm{ker}(dF_p)`$. Then there exists a smooth function $`s:U_p`$ such that the derivative of $`s`$ in the direction of $`o_p`$ is positive. (First construct $`s`$ along $`\mathrm{\Sigma }(F_p)`$ near $`\mathrm{\Sigma }^{1,1}(F_p)`$ then extend it to a tubular neighborhood of $`\mathrm{\Sigma }(F_p)`$.) The map $`F_p^{}=(F_p,s):U_p^6`$ is an immersion. If $`U_p`$ is chosen sufficiently small then we can even suppose that $`F_p^{}`$ is an embedding. Denote by $`e_6`$ the sixth coordinate direction in $`^6`$ and let $`\nu _6:U_pT^6`$ denote the vector field along $`F_p^{}`$ defined by the formula $`\nu _6(x)=e_6T_{F_p^{}(x)}^6`$ for $`xU_p`$. Projecting $`\nu _6`$ into the normal bundle of $`F_p^{}`$ we obtain a normal field $`\mu _6`$ along $`F_p^{}`$ that vanishes exactly at the points of $`\mathrm{\Sigma }(F_p)`$. Perturb $`F_p^{}`$ in the direction of $`\mu _6`$ to obtain an embedding $`F_p^{\prime \prime }`$. Then orient $`\mathrm{\Sigma }(F_p)`$ as the intersection of $`F_p^{}`$ and $`F_p^{\prime \prime }`$ in $`^6`$. Here $`^6`$ is considered with its standard orientation. This orientation of $`\mathrm{\Sigma }(F_p)`$ does not depend on the choice of the function $`s`$, since if $`s_1`$ and $`s_2`$ are two such functions then for $`0t1`$ the convex combination $`(1t)s_1+ts_2`$ also satisfies the conditions for $`s`$.
If we reverse the orientation of $`\mathrm{ker}(dF_p)`$, i.e. if we orient it by $`o_p`$, then we can choose $`s`$ instead of $`s`$. Thus we obtain the embedding $`(F_p,s)`$, which is the reflection of $`F_p^{}`$ in the hyperplane $`^5`$. Denote this reflection by $`R:^6^6`$ (i.e., $`R(x_1,\mathrm{},x_5,x_6)=(x_1,\mathrm{},x_5,x_6)`$). Then $`(F_p,s)=RF_p^{}`$. The vector field $`dR\nu _6`$ along $`RF_p^{}`$ points in the direction $`e_6`$ and $`RF_p^{\prime \prime }`$ is the perturbation of $`RF_p^{}`$ in the direction of $`dR\mu _6`$. But in this case we should perturb $`RF_p^{}`$ in the direction of $`(dR\mu _6)`$. We obtain the same orientation if we look at the intersection $`(RF_p^{\prime \prime })(RF_p^{})`$ instead. Since the intersection is 2-dimensional and $`U_p`$ is 4-dimensional we get that $`(RF_p^{\prime \prime })(RF_p^{})=(RF_p^{})(RF_p^{\prime \prime })`$ in the oriented sense. The orientations of $`\mathrm{\Sigma }(F_p)`$ defined by the intersections $`F_p^{}F_p^{\prime \prime }`$ and $`(RF_p^{})(RF_p^{\prime \prime })`$ are opposite. This can be seen from the following argument: For $`0t1`$ denote by $`R_t`$ the rotation of the hyperplane $`^6`$ in $`^7`$ around $`^5`$ by the angle $`\pi t`$. The orientation of $`(R_tF_p^{})(R_tF_p^{\prime \prime })`$ in $`R_t(^6)`$ changes continuously as $`t`$ goes from $`0`$ to $`1`$. The orientations of the hyperplanes $`R_1(^6)`$ and $`^6`$ are opposite, thus the reflection $`R`$ changes the orientation of the intersection $`F_p^{}F_p^{\prime \prime }`$.
So we have defined an isomorphism $`\mathrm{\Psi }_p`$ between $`\mathrm{ker}(dF_p)`$ and $`det(T\mathrm{\Sigma }(F_p))`$ for every $`pW^4`$ in a compatible way (i.e., $`\mathrm{\Psi }_p|(U_pU_q)=\mathrm{\Psi }_q|(U_pU_q)`$ for $`p,qW^4`$). These local isomorphisms define a global isomorphism $`\mathrm{\Psi }`$ between $`\mathrm{ker}(dF)`$ and $`det(T\mathrm{\Sigma }(F))`$.
## 5. The invariant
In this section we will give a geometric formula for the 3-dimensional obstruction to the existence of a regular homotopy between two immersions of $`M^3`$ into $`^5`$. This generalizes the results of to immersions with non-trivial normal bundle.
###### Definition 5.1.
Let $`W^4`$ be a compact oriented manifold with boundary $`M^3`$ and $`F:W^4^5`$ a generic map such that $`f=F|M^3`$ is an immersion. Recall that $`\mathrm{\Sigma }(F)`$ denotes the set of singular points of $`F`$. Let us denote by $`C(F)M^3`$ the 1-dimensional submanifold $`\mathrm{\Sigma }(F)`$. Choose a trivialization $`\tau `$ of $`\mathrm{ker}(dF)|C(F)`$ so that it points into the interior of $`W^4`$. This is possible since $`f`$ is non-singular (and so $`\mathrm{ker}(dF)`$ never lies in $`TM^3`$). Then $`\tau `$ is normal to $`\mathrm{\Sigma }(F)`$ because $`F`$ is generic and thus $`\mathrm{\Sigma }^{1,1}(F)C(F)=\mathrm{}`$. So if we project $`\tau `$ into $`TM`$ along $`\mathrm{\Sigma }(F)`$ we obtain a nowhere vanishing normal field $`\nu (F)`$ in $`\nu (C(F)M^3)`$.
Let $`U`$ denote a small collar neighborhood of $`C(F)`$ in $`\mathrm{\Sigma }(F)`$. Then clearly $`U`$ is orientable. Using the isomorphism $`\mathrm{\Psi }`$ of Definition 4.4 the trivialization $`\tau `$ of $`\mathrm{ker}(dF)`$ induces an orientation of $`U`$. Thus $`C(F)U`$ is also oriented. So we have assigned a pair $`(C(F),\nu (F))`$ to $`F`$ as in Notation 3.1. Let $`r(F)=r(C(F),\nu (F))`$.
###### Notation 5.2.
For $`\chi H^2(M^3;)`$ let us denote by $`\text{Imm}(M^3,^5)_\chi `$ the space of immersions with normal Euler class $`\chi `$.
Fix a cohomology class $`\chi H^2(M^3;)`$. Our aim is to define an invariant $`i:\pi _0\left(\text{Imm}(M^3,^5)_\chi \right)_{2d(\chi )}`$.
###### Proposition 5.3.
Let $`f\text{Imm}(M^3,^5)_\chi `$ and let $`F:W^4^5`$ be a generic map such that $`F=f`$. Then $`[C(F)]=D\chi `$.
###### Proof.
Let $`\kappa `$ denote an inner normal field of $`M^3`$ in $`W^4`$ that extends $`\tau `$ (see Definition 5.1). Then $`dF\kappa `$ is a vector field along $`f`$ that is tangent to $`f`$ exactly at the points of $`C(F)=\mathrm{\Sigma }(F)`$. (If $`pC(F)`$ then the rank of $`(dF)_p`$ is $`3`$, moreover $`dF|(T_pM^3)=df`$ is non-degenerate. Thus $`dF(\kappa _p)df(T_pM^3)`$.) So if we project $`dF\kappa `$ into the normal bundle of $`f`$ we obtain a normal field of $`f`$ that vanishes along $`C(F)`$. To see that $`C(F)`$ represents the normal Euler class of $`f`$, we have to know that it is oriented suitably.
Using the notations of Definition 5.1 we choose a function $`s:W^4`$ such that the derivative of $`s`$ in the direction of $`\kappa `$ (and thus $`\tau `$) is positive and $`s|M^30`$. Then there exists a collar neighborhood $`V`$ of $`M^3`$ in $`W^4`$ such that $`F^{}=(F,s)|V`$ is an immersion. Denote by $`\nu _F^{}`$ the normal bundle of $`F^{}`$ in $`^6`$ and by $`\nu _f`$ the normal bundle of $`f`$ in $`^5`$. Then $`\nu _F^{}|M^3=\nu _f`$ as oriented bundles, since $`s`$ is increasing along $`\kappa `$ (here $`^5`$ and $`^6`$ are considered with their standard orientations). By Definition 4.4 the surface of singular points $`U=\mathrm{\Sigma }(F|V)`$ is oriented as the self-intersection of $`F^{}`$ in $`^6`$, or more precisely, as the intersection of the zero section and a generic section of $`\nu _F^{}`$. Moreover, $`C(F)`$ is oriented as the boundary of $`U`$. Thus $`C(F)`$ is the self-intersection of the zero section of $`\nu _F^{}|M^3=\nu _f`$, so it is dual to the Euler class $`e(\nu _f)=\chi `$. (Here we used the naturality of the Euler class.) ∎
###### Definition 5.4.
Let $`f\text{Imm}(M^3,^5)_\chi `$ and let $`F:W^4^5`$ be generic such that $`F=f`$. Denote the algebraic number of cusps of $`F`$ by $`\mathrm{\#}\mathrm{\Sigma }^{1,1}(F)`$ (for the definition see ). Then let $`j(f)=3\sigma (W^4)3\alpha (M^3)+\mathrm{\#}\mathrm{\Sigma }^{1,1}(F)+r(F)_{2d(\chi )}`$.
Note that if $`\chi =0`$ and $`F`$ is an immersion near $`W^4`$ then $`r(F)=0`$. Thus in this case $`j(f)`$ agrees with the double of the invariant introduced in (see Definition 2.7).
###### Theorem 5.5.
j(f) is well defined, i.e., it does not depend on the choice of the generic map $`F`$. Moreover, if $`f_0`$ and $`f_1`$ are regularly homotopic then $`j(f_0)=j(f_1)`$.
###### Proof.
For $`i\{0,1\}`$ let $`F_i:W_i^4^5`$ be a generic map such that $`F_i=f_i`$. Choose a regular homotopy $`\{h_t:0t1\}`$ connecting $`f_0`$ and $`f_1`$. This defines an immersion $`H:M^3\times I^5\times I`$ by the formula $`H(x,t)=(h_t(x),t)`$. Also choose a closed collar neighborhood $`U_i`$ of $`M^3`$ in $`W_i^4`$ and a diffeomorphism $`d_i:U_iM^3\times [0,\epsilon ]`$ for $`i=0,1`$. Let $`p:M^3\times [0,\epsilon ][0,\epsilon ]`$ denote the projection onto the second factor. If $`\epsilon `$ (i.e., $`U_i`$) is sufficiently small then $`pd_i`$ is non-degenerate along $`\mathrm{ker}(dF_i)`$ for $`i=0,1`$ since $`\mathrm{ker}(dF_i)`$ never lies in $`TM^3`$. Let $`s_i`$ be an arbitrary smooth extension of $`pd_i`$ over $`W_i^4`$. Now let $`F_0^{}=(F_0,s_0):W_0^4^6`$ and $`F_1^{}=(F_1,s_1+1):W_1^4^6`$. Then $`F_i^{}`$ is an immersion on $`U_i`$. Notice that $`H|(M^3\times \{0\})=F_0^{}|(W_0^4)`$ and $`H|(M^3\times \{1\})=F_1^{}|(W_1^4)`$.
Denote by $`\kappa _i`$ the inner normal field of $`W_i^4`$ along $`M^3=W_i^4`$ and by $`v_6`$ the sixth coordinate direction in $`^6`$. Then the inner product $`dF_0^{}(\kappa _0),v_6<0`$ and $`dF_1^{}(\kappa _1),v_6>0`$. Furthermore, if $`\lambda _i`$ denotes the inner normal field of $`M^3\times I`$ along $`M^3\times \{i\}`$ for $`i=0,1`$ then $`dH(\lambda _0),v_6>0`$ and $`dH(\lambda _1),v_6<0`$. So $`dH(\lambda _i)`$ is homotopic to $`dF_i(\kappa _i)`$ in the space of vector fields normal to $`H|(M^3\times \{i\})`$. Using Smale’s lemma there exists a regular homotopy of $`H`$ fixed on the boundary $`M^3\times \{0,1\}`$ that induces the above homotopy of normal fields. Denote by $`H^{}`$ the result of this regular homotopy of $`H`$. Then $`F_0^{}`$, $`H^{}`$ and $`F_1^{}`$ fit together to a smooth map $`F^{}`$ of $`W^4=W_0^4(M^3\times I)W_1^4`$ into $`^6`$ that is an immersion on $`M^3\times I`$. Let $`\pi :^6^5`$ denote the projection map. Then by a small perturbation of $`H^{}`$ we can achieve that $`F=\pi F^{}`$ is generic.
Since $`G=F|(M^3\times I)=\pi H^{}`$ is prim, the singular surface $`\mathrm{\Sigma }(G)`$ is oriented and a trivialization $`\tau `$ of $`\mathrm{ker}(dG)`$ is given. If we project $`\tau `$ into $`\nu (\mathrm{\Sigma }(G)M^3\times I)`$ we obtain a normal field $`\nu `$ along $`\mathrm{\Sigma }(G)`$ that vanishes exactly at the cusps of $`G`$, i.e., where $`\tau `$ is tangent to $`\mathrm{\Sigma }(G)`$. So $`\mathrm{\#}\mathrm{\Sigma }^{1,1}(G)`$ is equal to the algebraic number of zeroes of $`\nu `$, which in turn is congruent to $`\text{rd}((C(F_0),\nu (F_0)),(C(F_1),\nu (F_1)))=r(F_0)r(F_1)`$ modulo $`2d(\chi )`$ by Definition 3.2.
Now using the result of Szűcs that $`3\sigma (W^4)+\mathrm{\#}\mathrm{\Sigma }^{1,1}(F)=0`$ we get that
(5.1)
$$\begin{array}{c}\left(3\sigma (W_0^4)+\mathrm{\#}\mathrm{\Sigma }^{1,1}(F_0)+r(F_0)\right)+\left(3\sigma (W_1^4)+\mathrm{\#}\mathrm{\Sigma }^{1,1}(F_1)+r(F_1)\right)=\hfill \\ \hfill 3\sigma (W^4)+\mathrm{\#}\mathrm{\Sigma }^{1,1}(F_0GF_1)=0.\end{array}$$
In the special case $`f_0=f_1=f`$ this implies that $`j(f)`$ is well defined, and for $`f_0`$ and $`f_1`$ arbitrary (but regularly homotopic) we get that $`j`$ is a regular homotopy invariant. ∎
###### Proposition 5.6.
For any immersion $`f:M^3^5`$ the invariant $`j(f)`$ is always an even element of $`_{2d(\chi )}`$, i.e., it is mapped to $`0`$ by the epimorphism $`_{2d(\chi )}_2`$.
###### Proof.
Choose an immersion $`f_1\text{Imm}(M^3,^5)_0`$ and denote $`f`$ by $`f_0`$. Since in it is proved that $`j(f_1)`$ is always even ($`j(f_1)=2i(f_1)`$ for $`i`$ as in Definition 2.7) it is sufficient to prove that $`j(f_0)j(f_1)mod2`$. Choose a singular Seifert surface $`F_i:W_i^5`$ for $`f_i`$ ($`i=0,1`$) and let $`G:M^3\times I^5`$ be a generic map such that $`F_0GF_1`$ is a smooth map on $`W_1(M^3\times I)W_2`$. Then by equation 5.1 above it is sufficient to prove that $`r(F_0)r(F_1)\mathrm{\#}\mathrm{\Sigma }^{1,1}(G)mod2`$. Since $`f_1`$ has trivial normal bundle we may choose $`F_1`$ to be an immersion in a neighborhood of $`W_1^4`$. So $`G`$ is an immersion in a neighborhood of $`M^3\times \{1\}`$, moreover $`r(F_1)=0`$. The difference between the present situation and the proof of Theorem 5.5 is that now $`\mathrm{ker}(dG)`$ might be non-orientable. Using Definition 3.8 we get that $`r_2(F_0)r_2(F_1)r(F_0)r(F_1)mod2`$. Let $`\nu `$ denote a generic normal field along $`\mathrm{\Sigma }^1(G)`$ that extends both $`\nu (F_0)`$ and $`\nu (F_1)`$. By definition $`r_2(F_0)r_2(F_1)`$ equals the mod $`2`$ number of zeroes of $`\nu `$. Thus we only have to prove that $`\left|\nu ^1(0)\right|\mathrm{\#}\mathrm{\Sigma }^{1,1}(G)mod2`$.
From now on we will denote $`\mathrm{\Sigma }^1(G)`$ by $`K`$ and the line bundle $`\mathrm{ker}(dG)<T(M^3\times I)|K`$ by $`l`$. Then $`l`$ is tangent to $`K`$ exactly at the points of $`\mathrm{\Sigma }^{1,1}(G)`$. For $`\epsilon >0`$ sufficiently small let $`\stackrel{~}{K}`$ denote the sphere bundle $`S_\epsilon l`$. If $`\epsilon `$ is sufficiently small then the exponential map of $`M^3\times I`$ defines an immersion $`s:\stackrel{~}{K}M^3\times I`$ so that the double points of $`s`$ correspond exactly to the points of $`\mathrm{\Sigma }^{1,1}(G)`$. So we have to prove that $`\left|D_2(s)\right|\left|\nu ^1(0)\right|mod2`$. By Lemma 4.1 the surface $`\stackrel{~}{K}`$ is the orientation double cover of $`K`$, in particular $`\stackrel{~}{K}`$ is oriented and a sign can be given to each double point of $`s`$ (here we also use that $`dim(\stackrel{~}{K})`$ is even). The sign of a double point of $`s`$ is the opposite of the sign of the corresponding cusp of $`G`$ (since the sign of a cusp is defined as the self intersection of $`K`$). Thus $`\mathrm{\#}D_2(s)=\mathrm{\#}\mathrm{\Sigma }^{1,1}(G)`$. Let $`p:\stackrel{~}{K}K`$ denote the covering map. Then $`p^{}\nu _K\nu _s`$, thus $`p^{}\nu `$ defines a section $`\stackrel{~}{\nu }`$ of $`\nu _s`$. From the construction of $`\stackrel{~}{\nu }`$ it is clear that $`\mathrm{\#}\stackrel{~}{\nu }^1(0)=2\mathrm{\#}\nu ^1(0)`$.
If we perturb $`s`$ in the direction of $`\stackrel{~}{\nu }`$ we get a self intersection point of $`s`$ for each element of $`\stackrel{~}{\nu }^1(0)`$ and two self intersection points for each double point of $`s`$. Thus
$$s(s+\epsilon \stackrel{~}{\nu })=\mathrm{\#}\stackrel{~}{\nu }^1(0)+2D_2(s)=2(\mathrm{\#}\nu ^1(0)\mathrm{\#}\mathrm{\Sigma }^{1,1}(G)).$$
So we only have to show that the left hand side is divisible by $`4`$.
From now on we will work in a fixed tubular neighborhood $`T`$ of $`C(F_0)M^3`$. Note that $`(s,\stackrel{~}{\nu })=(C(F_0)+\epsilon \nu (F_0),\nu (F_0))(C(F_0)\epsilon \nu (F_0),\nu (F_0))T\times \{0\}`$. Let us denote by $`C`$ the one-dimensional submanifold $`(C(F_0)+\epsilon \nu (F_0))(C(F_0)\epsilon \nu (F_0))`$ of $`M^3`$. Then $`s(s+\epsilon \stackrel{~}{\nu })=r(C,\nu (F_0))`$ since $`CC(F_0)C(F_0)`$ is null homologous in $`M^3`$. We define an embedding $`e:C(F_0)\times [\epsilon ,\epsilon ]T`$ by the formula $`e(x,t)=x+t\nu (F_0)`$. Then $`E=\text{Im}(e)`$ is a 2-dimensional oriented submanifold of $`T`$ with boundary $`C`$. Thus $`r(C,\nu (F_0))=E(C+\nu (F_0))`$ where the right hand side is considered to be a generic intersection (each fiber of $`E`$ is parallel to $`\nu (F_0)`$). Let $`n`$ be a small non-zero vector field along $`C(F_0)`$ orthogonal to $`\nu (F_0)`$. Then $`E(C+\nu (F_0)+n)=\mathrm{}`$ (this can be verified by inspecting each fiber of $`T`$), thus $`r(C,\nu (F_0))=0`$. So we get that $`\mathrm{\#}\nu ^1(0)=\mathrm{\#}\mathrm{\Sigma }^{1,1}(G)`$, not just a mod 2 congruence. ∎
###### Remark 5.7.
A small improvement on the proof of Proposition 5.6 yields an interesting result: Let $`G:M^3\times [0,1]^5`$ be a generic map connecting the immersions $`f_0`$ and $`f_1`$. Let $`K`$ denote the singular surface of $`G`$ and let $`\nu _i`$ be a trivialization of $`\mathrm{ker}(dG)|(M^3\times \{i\})`$ for $`i=0,1`$. Then $`\mathrm{\#}\mathrm{\Sigma }^{1,1}(G)`$ is equal to the relative twisted normal Euler class $`e(\nu _K;\nu _0,\nu _1)`$.
Note that if $`K=C_0C_1`$ then by definition
$$\text{rd}_2((C_0,\nu _0),(C_1,\nu _1))e(\nu _K;\nu _0,\nu _1)mod2.$$
If in particular $`K`$ is an oriented cobordism between $`C_0`$ and $`C_1`$ then
$$\text{rd}((C_0,\nu _0),(C_1,\nu _1))=e(\nu _K;\nu _0,\nu _1)=\mathrm{\#}\mathrm{\Sigma }^{1,1}(G);$$
here $`C_i`$ is oriented by $`\nu _i`$ using the isomorphism $`\mathrm{\Psi }`$ (see Definition 4.4).
Thus $`j`$ may take only $`d(\chi )`$ different values if $`d(\chi )>0`$. (Since $`j`$ is additive if a connected sum is taken with an immersion of a sphere (Lemma 6.1) and $`j(g)`$ can be any even number for $`g:S^3^5`$ it follows that $`j`$ is an epimorphism onto $`2_{2d(\chi )}`$.) But Theorem 2.1 implies that there are exactly $`2d(\chi )`$ regular homotopy classes with normal Euler class $`\chi `$. So $`j`$ describes the regular homotopy class of $`f`$ only up to a $`2:1`$ ambiguity. To resolve this problem we will lift the invariant $`j_{2d(\chi )}`$ to an invariant $`I_{4d(\chi )}`$. It follows from Proposition 5.6 that $`I`$ is always an even element, thus it defines an invariant $`i_{2d(\chi )}`$ by the embedding $`_{2d(\chi )}_{4d(\chi )}`$.
###### Notation 5.8.
Let $`f\text{Imm}(M^3,^5)_\chi `$ and let $`F:W^4^5`$ be a singular Seifert surface for $`f`$. If $`\kappa `$ denotes the inner normal field along $`W^4`$ then let $`\overline{w}(F)\mathrm{\Gamma }(\nu _f)`$ be the projection of $`dF(\kappa )`$ into $`\nu _f`$. If $`\overline{t}\mathrm{\Gamma }(\epsilon _M^1)`$ denotes a trivialization of $`\epsilon _M^1`$ then we can consider $`\overline{t}`$ and $`\overline{w}(F)`$ to be sections of $`\nu _f\epsilon _M^1`$. Let $`\overline{v}(F)=\overline{w}(F)+\overline{t}\mathrm{\Gamma }(\nu _f\epsilon _M^1)`$.
From now on we will fix a spin structure $`s_M\text{Spin}(M^3)`$. If we consider $`^5`$ with its unique spin structure then for every immersion $`f:M^3^5`$ a spin structure $`s(f)`$ is induced on $`\nu _f`$ by $`s_M`$. Then $`s(f)`$ is equivalent to a trivialization $`\tau (f):\epsilon _M^3|\text{sk}_2(M)\nu _f\epsilon _M^1|\text{sk}_2(M^3)`$ up to homotopy. Since $`\pi _2(SO(3))=0`$ the trivialization $`\tau (f)`$ extends to an isomorphism $`\tau (f):\epsilon _M^3\nu _f\epsilon _M^1`$, but this extension is not unique because $`\pi _3(SO(3))0`$.
###### Definition 5.9.
Using the above notations let $`t(f),v(F)\mathrm{\Gamma }(\epsilon _M^3)`$ be defined by the formulas $`t(f)=\tau (f)^1\overline{t}`$ and $`v(F)=\tau (f)^1\overline{v}(F)`$. Denote by $`a(F)`$ the quadruple $`(C(F),\nu (F),t(f),v(F))N(M^3,\chi )`$. Then define $`R(F)_{4d(\chi )}`$ to be $`R(a(F))`$. Since the homotopy class of the map $`(t(f),v(F))|\text{sk}_2(M^3):\text{sk}_2(M^3)V_2(^3)`$ is independent of the choice of the extension of $`\tau (f)`$ to $`\epsilon _M^3`$ Proposition 3.14 implies that $`R(F)`$ is also independent of $`\tau (f)`$ and depends only on $`s_M`$.
###### Remark 5.10.
Proposition 3.15 implies that $`r(F)R(F)mod2d(\chi )`$.
Now we can finally define a complete regular homotopy invariant.
###### Definition 5.11.
For $`f\text{Imm}(M^3,^5)_\chi `$ and a singular Seifert surface $`F`$ let $`I(f)_{4d(\chi )}`$ be defined as $`3\sigma (W^4)3\alpha (M^3)+\mathrm{\#}\mathrm{\Sigma }^{1,1}(F)+R(F)`$. (Recall that we have fixed a spin structure $`s_M`$ on $`M^3`$ for the definition of $`R(F)`$.) Remark 5.10 above implies that $`j(f)I(F)mod2d(\chi )`$. Thus by Proposition 5.6 we get that $`I(F)`$ is always an even element of $`_{4d(\chi )}`$. Let us denote by $`\frac{1}{2}`$ the isomorphism from $`2_{4d(\chi )}`$ to $`_{2d(\chi )}`$. Then let $`i(F)=\frac{1}{2}I(F)`$.
Clearly $`j(f)=2i(f)`$ for every $`f\text{Imm}(M^3,^5)_\chi `$.
###### Theorem 5.12.
I(f) is well defined, i.e., it does not depend on the choice of the generic map $`F`$. Moreover, if $`f_0`$ and $`f_1`$ are regularly homotopic then $`I(f_0)=I(f_1)`$.
###### Proof.
Using the notations of the proof of Theorem 5.5 we only have to show that the surface $`K=\mathrm{\Sigma }(G)M^3\times I`$ satisfies Definition 3.11. I.e., there exist generic sections $`t`$ and $`v`$ of $`\epsilon _{M\times I}^3`$ that extend $`t(f_i)`$ and $`v(F_i)`$ for $`i=0,1`$ and are linearly dependent exactly over $`K`$. The regular homotopy between $`f_0`$ and $`f_1`$ defines the immersion $`H:M^3\times I^5\times I`$. For $`i\{0,1\}`$ there is a canonic isomorphism $`\phi _i:\nu _H|(M^3\times \{i\})\nu _{f_i}`$. Let $`\overline{w}\mathrm{\Gamma }(\nu _H)`$ denote the projection of the sixth coordinate vector $`v_6^6`$ into $`\nu _H`$. Then $`\phi _i(\overline{w}|M^3\times \{i\})=\overline{w}(F_i)`$. Moreover, $`K=\overline{w}^1(0)`$ and the orientation of $`K`$ is defined as the self intersection of $`H`$ if perturbed in the direction of $`\overline{w}`$. Define $`\overline{t}`$ to be a trivialization of the $`\epsilon _{M\times I}^1`$ component of the bundle $`\nu _H\epsilon _{M\times I}^1`$ and let $`\overline{v}=\overline{w}+\overline{t}`$. Then $`\overline{t}`$ and $`\overline{v}`$ are linearly dependent exactly at the points of $`K`$. Note that $`t(f_i)=\tau (f_i)^1(\phi _i\text{id}_{\epsilon ^1})(\overline{t}|M^3\times \{i\})`$ and $`v(F_i)=\tau (f_i)^1(\phi _i\text{id}_{\epsilon ^1})(\overline{v}|M^3\times \{i\})`$. Thus we only have to define a trivialization $`\tau :\epsilon _{M\times I}^3\nu _H\epsilon _{M\times I}^1`$ such that
(5.2)
$$\tau |(M^3\times \{i\})=\phi _i^1\tau (f_i)\text{for}i=0,1.$$
The spin structure $`s_M\text{Spin}(M^3)`$ and the unique spin structure on $`I`$ define a spin structure on $`M^3\times I`$. Together with the unique spin structure of $`^6`$ we get a spin structure $`s_H`$ on $`\nu _H`$. When $`s_H`$ is restricted to $`M^3\times \{i\}`$ we get back the spin structure $`s_M`$. Thus $`s_H`$ defines a trivialization
$$\tau _H:\epsilon _{M\times I}^3|\text{sk}_2(M^3\times I)(\nu _H\epsilon _{M\times I}^1)|\text{sk}_2(M^3\times I)$$
satisfying equation 5.2 over the 2-skeleton of $`M^3\times \{0,1\}`$. Note that the trivialization $`\tau (f_i)`$ is only well defined over $`\text{sk}_2(M^3)`$ and that we can choose an arbitrary extension over $`M^3`$ in order to define the rotation difference. Thus we only have to extend $`\tau _H`$ to a trivialization $`\tau `$ of $`\nu _H\epsilon _{M\times I}^1`$ and then define $`\tau (f_i)`$ by formula 5.2.
First we extend $`\tau _H`$ to $`\text{sk}_3(M^3\times I)\text{sk}_3(M^3\times \{1\})`$. This is possible since the obstruction to extending the trivialization over a 3-simplex from its boundary lies in $`\pi _2(SO(3))=0`$. If $`\sigma ^3`$ is a 3-simplex of $`M^3`$ then we can extend $`\tau _H`$ to $`\sigma ^3\times I`$ since it is given only on $`(\sigma ^3\times I)(\sigma ^3\times \{1\})`$. Thus we have obtained the required extension $`\tau `$ of $`\tau _H`$. ∎
## 6. Connected sums and completeness of the invariant $`i`$
###### Lemma 6.1.
If $`f\text{Imm}(M^3,^5)_\chi `$ and $`g\text{Imm}(S^3,^5)`$ then $`c(f\mathrm{\#}g)=c(f)`$, in particular $`e(\nu _{f\mathrm{\#}g})=e(\nu _f)`$. Moreover
$$i(f\mathrm{\#}g)=i(f)+\left(i(g)mod2d(\chi )\right)_{2d(\chi )}.$$
###### Proof.
Since $`c(f)`$ describes the regular homotopy class of $`f|(M_{}^3)`$ it is trivial that $`c(f\mathrm{\#}g)=c(f)`$. Let $`F`$ be a singular Seifert surface of $`f`$ and $`G`$ of $`g`$ such that $`G`$ is an immersion near the boundary. Then the result follows by inspecting the boundary connected sum $`F\mathrm{}G`$ and the fact that $`C(G)=\mathrm{}`$. ∎
###### Theorem 6.2.
Suppose that the immersions $`f_0,f_1\text{Imm}(M^3,^5)`$ are regularly homotopic on $`M^3D`$, where $`DM^3`$ is diffeomorphic to the closed disc $`D^3`$ (i.e., $`c(f_0)=c(f_1)`$). Then $`i(f_0)=i(f_1)`$ implies that $`f_0`$ is regularly homotopic to $`f_1`$.
###### Proof.
The proof consists of two cases according to the value of $`d(\chi )`$.
If $`d(\chi )>0`$ then $`i`$ takes values in $`_{2d(\chi )}`$ which is a finite group. Theorem 2.1 implies that there are exactly $`2d(\chi )`$ regular homotopy classes with a fixed Wu invariant $`c`$. Thus we only have to show that the invariant $`i`$ restricted to immersions with Wu invariant $`c`$ is an epimorphism onto $`_{2d(\chi )}`$. For this end choose an immersion $`f\text{Imm}(M^3,^5)`$ such that $`c(f)=c`$. In it is shown that $`i:\text{Imm}[S^3,^5]`$ is a bijection. Thus Lemma 6.1 implies that $`c(f\mathrm{\#}g)=c(f)=c`$ for every $`g\text{Imm}(S^3,^5)`$, moreover $`i:\{f\mathrm{\#}g:g\text{Imm}(S^3,^5)\}_{2d(\chi )}`$ is surjective.
If $`d(\chi )=0`$ then $`i`$ maps into $``$. Using Smale’s lemma we can suppose that $`f_0|(M^3D)=f_1|(M^3D)`$. The normal bundles of $`f_0|D`$ and $`f_1|D`$ in $`^5`$ are trivial, choose a trivialization for both of them. Let $`\tau _0`$ be a non-zero normal field along $`f_0|D`$. Then $`\tau _0|D`$ considered in the trivialization of the normal bundle of $`f_1|D`$ is a map $`(\tau _0|D):DS^1`$. Since $`D`$ is homeomorphic to $`S^2`$ and $`\pi _2(S^1)=0`$ the normal field $`\tau _0|D`$ can be extended to a normal field $`\tau _1`$ of $`f_1|D`$. Thus $`\tau _i`$ is a normal field of $`f_i|D`$ for $`i=0,1`$ and $`\tau _0|D=\tau _1|D`$.
Next choose an oriented compact manifold $`W_0^4`$ with boundary $`M^3`$. We push $`D`$ into the interior of $`W_0^4`$ fixing the boundary $`D`$ to obtain a 3-disc $`D_1W_0^4`$ so that $`D=D_1`$ and $`M_1^3=(M^3D)D_1`$ is a smooth submanifold of $`W_0^4`$. If we throw out the domain bounded by $`D`$ and $`D_1`$ in $`W_0^4`$ we obtain a 4-dimensional submanifold $`W_1^4`$ of $`W_0^4`$ with boundary $`M_1^3`$. Clearly $`W_0^4`$ is diffeomorphic to $`W_1^4`$.
We can choose a generic map $`F_0:W_0^4^5`$ with the following three properties:
1. $`F_0|M^3=f_0`$ and $`F_0|M_1^3=f_1`$ (where $`M_1^3`$ is identified with $`M^3`$ by a diffeomorphism keeping $`M^3D`$ fixed).
2. $`F_0`$ is an immersion in a neighborhood of $`D`$ and $`D_1`$.
3. If $`\kappa _0`$ denotes the inner normal field of $`D`$ in $`W_0^4`$ and $`\kappa _1`$ denotes the inner normal field of $`D_1`$ in $`W_1^4`$ then $`dF_0\kappa _0=\tau _0`$ and $`dF_1\kappa _1=\tau _1`$.
Let $`F_1=F_0|W_1^4`$. Then (2) implies that $`C(F_0)=C(F_1)M^3D`$, moreover $`\nu (F_0)=\nu (F_1)`$. In particular, the normal Euler class of $`f_0`$ and $`f_1`$ coincide. Thus $`R(F_0)=R(F_1)`$. Since $`\sigma (W_0^4)=\sigma (W_1^4)`$, we get that
$$0=i(f_0)i(f_1)=\mathrm{\#}\mathrm{\Sigma }^{1,1}(F_0|(W_0^4W_1^4)).$$
Choose diffeomorphisms $`d_0:S_+^3D`$ and $`d_1:S_+^3D_1`$, where $`S_+^3`$ denotes the northern hemisphere of $`S^3`$. Then the immersion $`F_0d_i`$ can be extended to an immersion $`f_i^{}:S^3^5`$ for $`i=0,1`$ so that $`f_0^{}|S_{}^3=f_1^{}|S_{}^3`$. (This is possible since $`j^1(f_0)|D=j^1(f_1)|D.`$) Now repeating the same argument as above for $`f_0^{}`$ and $`f_1^{}`$, we obtain that
$$i(f_0^{})i(f_1^{})=\mathrm{\#}\mathrm{\Sigma }^{1,1}(F_0|(W_0^4W_1^4)).$$
(Note that $`\tau _0`$ and $`\tau _1`$ have a common extension over $`S_{}^3`$.) Thus $`i(f_0^{})i(f_1^{})=0`$, so using we get that $`f_0^{}`$ and $`f_1^{}`$ are regularly homotopic. But this implies that there exists a regular homotopy between $`f_0^{}`$ and $`f_1^{}`$ that is fixed on $`S_{}^3`$ (see , Lemma 3.33). So $`f_0|D`$ and $`f_1|D`$ are regularly homotopic keeping the 1-jets on the boundary fixed, which completes the proof that $`f_0`$ and $`f_1`$ are regularly homotopic. ∎
###### Corollary 6.3.
The map
$$(c,i):\text{Imm}[M^3,^5]\underset{cH^2(M^3;)}{}_{4d(c)}$$
is a bijection.
We get more structure on the set of regular homotopy classes of immersions of oriented 3-manifolds into $`^5`$ if we endow it with the connected sum operation. Let us introduce the notation
$$I(3,5)=\{[f]:[f]\text{Imm}[M^3,^5]\text{for}M^3\text{oriented}\}.$$
Then $`(I(3,5),\mathrm{\#})`$ is a semigroup whose structure is described in the following theorem.
###### Theorem 6.4.
Let $`M_1^3`$ and $`M_2^3`$ be oriented 3-manifolds. Then
(6.1)
$$H^2(M_1^3\mathrm{\#}M_2^3;)H^2(M_1^3;)H^2(M_2^3;).$$
If $`f_i\text{Imm}(M_i^3,^5)`$ for $`i=1,2`$ then
(6.2)
$$c(f_1\mathrm{\#}f_2)=c(f_1)c(f_2)H^2(M_1^3\mathrm{\#}M_2^3;).$$
Moreover, if $`\chi _i`$ denotes the normal euler class of $`f_i`$ and $`\chi `$ the normal euler class of $`f_1\mathrm{\#}f_2`$ then $`d(\chi )=\mathrm{gcd}(d(\chi _1),d(\chi _2))`$. Finally,
(6.3)
$$i(f_1\mathrm{\#}f_2)=(i(f_1)mod2d(\chi ))+(i(f_2)mod2d(\chi ))_{2d(\chi )},$$
where $`i(f_i)_{2d(\chi _i)}`$ for $`i=1,2`$.
###### Proof.
Equation 6.1 follows from the fact that $`H^2(M_i^3;)H^2(M_i^3D^3;)`$ (see the long exact sequence of the pair $`(M_i^3,M_i^3D^3)`$) and the Mayer-Vietoris exact sequence for $`M_1^3\mathrm{\#}M_2^3=(M_1^3D^3)(M_2^3D^3)`$.
Equation 6.2 can be seen from the description of $`c(f_i)`$ as the regular homotopy class of $`f_i|\text{sk}_2(M^3)`$. Since $`\chi =\chi _1\chi _2`$ the statement about $`d(\chi )`$ is trivial.
Finally, equation 6.3 is obtained by taking the boundary connected sum $`F_1\mathrm{}F_2`$ of singular Seifert surfaces $`F_1`$ and $`F_2`$ for $`f_1`$, respectively $`f_2`$. ∎
## 7. Immersions of $`M^3`$ into $`^6`$ with a normal field
Let $`\text{Imm}_1(M^3,^6)`$ denote the space of immersions of $`M^3`$ into $`^6`$ with a normal field $`\nu `$. Moreover, let $`\text{Imm}_1[M^3,^6]=\pi _0(\text{Imm}_1(M^3,^6))`$ be the set of regular homotopy classes of such immersions with normal fields. If we fix a trivialization of $`TM^3`$ then Hirsch’s theorem implies that the natural map $`\text{Imm}_1(M^3,^6)C(M^3,V_4(^6))`$ is a weak homotopy equivalence.
For $`f\text{Imm}(M^3,^5)`$ let $`\iota (f)\text{Imm}_1(M^3,^6)`$ be the immersion $`f`$ with the constant normal field defined by the sixth coordinate vector in $`^6`$. Thus $`\iota `$ is an embedding of $`\text{Imm}(M^3,^5)`$ into $`\text{Imm}_1(M^3,^6)`$. As a special case of Hirsch’s compression theorem we have the following proposition.
###### Proposition 7.1.
$`\iota _{}:\text{Imm}[M^3,^5]\text{Imm}_1[M^3,^6]`$ is a bijection.
###### Proof.
The embedding $`^5^6`$ induces an embedding $`V_3(^5)V_4(^6)`$ and thus a map $`\psi :[M^3,V_3(^5)][M^3,V_4(^6)]`$ that makes the following diagram commutative.
$$\begin{array}{ccc}\text{Imm}[M^3,^5]& \stackrel{\iota _{}}{}& \text{Imm}_1[M^3,^6]\\ & & & & \\ [M^3,V_3(^5)]& \stackrel{\psi }{}& [M^3,V_4(^6)].\end{array}$$
By Hirsch’s theorem the vertical arrows are bijections, thus it is sufficient to prove that $`\psi `$ is also a bijection. To see this consider the fibration $`V_3(^5)V_4(^6)S^5`$. Then from the homotopy exact sequence of this fibration we get that the homomorphism $`\pi _i(V_3(^5))\pi _i(V_4(^6))`$ is an isomorphism for $`i3`$ and this implies that $`\psi `$ is a bijection. ∎
The natural forgetful map $`\phi :\text{Imm}_1(M^3,^6)\text{Imm}(M^3,^6)`$ is a Serre fibration.
###### Proposition 7.2.
For any immersion $`f:M^3^6`$ the normal bundle $`\nu _f`$ is trivial.
###### Proof.
Since $`M^3`$ is spin the normal bundle $`\nu _f`$ is also spin, thus it is trivial over the 2-skeleton of $`M^3`$. Such a trivialization can be extended to the 3-simplices of $`M^3`$ because $`\pi _2(SO(3))=0`$. ∎
So $`\phi `$ is surjective and the fiber of $`\phi `$ is homotopy equivalent to $`\mathrm{\Gamma }(\nu _f)=C(M^3,S^2)`$. Thus the end of the homotopy exact sequence of $`\phi `$ looks like as follows:
$$\pi _1(\text{Imm}(M^3,^6))[M^3,S^2]\text{Imm}_1[M^3,^6]\stackrel{\phi _{}}{}\text{Imm}[M^3,^6]0.$$
By Hirsch’s theorem there is a bijection $`\text{Imm}[M^3,^6][M^3,V_3(^6)]`$. Since $`V_3(^6)`$ is 2-connected and $`\pi _3(V_3(^6))_2`$ we get from obstruction theory that $`[M^3,V_3(^6)]H^3(M^3;_2)_2`$. It is well known that for $`f\text{Imm}(M^3,^6)`$ the regular homotopy class of $`f`$ is determined by the number of its double points $`D(f)`$ modulo 2. This gives a geometric interpretation of the map $`\phi _{}`$: for $`(f,\nu )\text{Imm}_1(M^3,^6)`$ the regular homotopy invariant $`\phi _{}(f,\nu )`$ is equal to $`D(f)`$ modulo 2.
How can we determine the value of $`\phi _{}\iota _{}(g)`$ for a generic $`g\text{Imm}(M^3,^5)\mathrm{?}`$ This question was answered in , let us recall that result now. The self-intersection set $`A(g)`$ of $`g`$ is a closed 1-dimensional submanifold of $`M^3`$ and $`g(A(g))`$ is also a closed 1-dimensional submanifold of $`^5`$. We say that a component $`C`$ of $`g(A(g))`$ is non-trivial if the double cover $`g|g^1(C):g^1(C)C`$ is non-trivial, i.e., if $`g^1(C)`$ is connected. Let the number of non-trivial components be denoted by $`\delta (g)`$. In Szűcs proved the following.
###### Theorem 7.3.
Suppose that $`f:M^3^6`$ is a generic immersion and that $`\pi :^6^5`$ is a projection such that $`g=\pi f`$ is also a generic immersion. Then
$$D(f)\delta (g)mod2.$$
Note that the immersions $`f`$ and $`g`$ above are regularly homotopic in $`^6`$. Thus for any generic $`g`$ we have that $`\phi _{}\iota _{}=\delta `$.
Now we are going to determine the group $`\pi _1(\text{Imm}(M^3,^6))`$. Using Hirsch’s theorem we get that it is isomorphic to $`\pi _1(C(M^3,V_3(^6)))=[SM^3,V_3(^6)]`$. Here $`SM^3`$ denotes the suspension of $`M^3`$ and is a 4-dimensional CW complex. The space $`V_3(^6)`$ is 2-connected and $`\pi _3(V_3(^6))_2`$. Moreover, $`\pi _4(V_3(^6))0`$. This can be seen as follows: From the homotopy exact sequence of the fibration $`V_3(^6)V_4(^7)S^6`$ we get that $`\pi _4(V_3(^6))\pi _4(V_4(^7))`$. It was shown by Paechter that for $`k4`$ the isomorphism $`\pi _k(V_k(^{2k1}))0`$ holds if $`k0mod4`$. Thus obstruction theory yields that $`[SM^3,V_3(^6)]H^3(SM^3;_2)H^2(M^3;_2)`$.
Putting together the above results we obtain the following theorem.
###### Theorem 7.4.
The following sequence is exact:
$$H^2(M^3;_2)[M^3,S^2]\text{Imm}[M^3,^5]\stackrel{𝛿}{}_20.$$
###### Remark 7.5.
If we fix a trivialization of $`TM^3`$ then non-zero vector fields (or equivalently, oriented 2-plane fields) on $`M^3`$ correspond to maps $`M^3S^2`$. Thus the set of homotopy classes of oriented 2-plane fields on $`M^3`$ is equal to $`[M^3,S^2]`$ which was determined in Remark 3.3. A geometric classification of such oriented 2-plane fields, avoiding the use of a trivialization of $`TM^3`$, was carried out by Gompf (see , section 4). A complete set of homotopy invariants, similar to those introduced in our present paper, were obtained in . Gompf’s result and the regular homotopy classification of immersions of $`M^3`$ into $`^5`$ are related by Theorem 7.4.
## Acknowledgements
I would like to take this opportunity to thank Professor András Szűcs for our long and helpful discussions and for reading earlier versions of this paper. |
warning/0506/math0506174.html | ar5iv | text | # A Characteristic Number of Hamiltonian Bundles over 𝑆²
## 1. Introduction
A loop $`\psi :S^1\text{Ham}(M,\omega )`$ in the group of Hamiltonian diffeomorphisms of a symplectic manifold $`(M^{2n},\omega )`$ can be considered as a clutching function of a Hamiltonian fibration $`E\stackrel{\pi }{}S^2`$ with fibre $`M`$. The total space $`E`$ supports the coupling class $`cH^2(E,)`$; this is the unique class such that $`c^{n+1}=0`$, and $`i_p^{}(c)`$ is the cohomology class of the symplectic structure on the fibre $`\pi ^1(p)`$, where $`i_p`$ is the inclusion of $`\pi ^1(p)`$ in $`E`$ . Furthermore one can consider on $`E`$ the first Chern class $`c_1(VTE)`$ of the vertical tangent bundle of $`E`$. These canonical cohomology classes on $`E`$ determine the characteristic number (see )
(1.1)
$$I_\psi =_Ec_1(VTE)c^n,$$
which depends only on the homotopy class of $`\psi `$. Since $`I`$ is an $``$-valued group homomorphism on $`\pi _1(\text{Ham}(M,\omega ))`$, the non vanishing of $`I`$ implies that the group $`\pi _1(\text{Ham}(M,\omega ))`$ is infinite. That is, $`I`$ can be used to detect the infinitude of the corresponding homotopy group. Furthermore $`I`$ calibrates the Hofer’s norm $`\nu `$ on $`\pi _1(\text{Ham}(M,\omega ))`$ in the sense that $`\nu (\psi )C|I_\psi |`$, for all $`\psi `$, where $`C`$ is a positive constant .
$`I`$ is a generalization of the mixed action-Maslov homomorphism introduced by Polterovich for monotone manifolds, that is, when $`[\omega ]=ac_1(TM)`$ and $`a>0`$. The value of this mixed action-Maslov homomorphism on a loop $`\psi `$ is, in many cases, easy to calculate, since it is a linear combination of the symplectic action around any orbit $`\{\psi _t(x_0)\}_t`$ and the Maslov index of the linearized flow $`(\psi _t)_{}`$ along this orbit. By contrast, $`I`$ is defined for Hamiltonian loops in general manifolds (non necessarily monotone), and its value is mostly not so easy to determine from the definition.
Our purpose in this note is to obtain an explicit expression for $`I_\psi `$, which can be used to calculate its value. More precisely, when the bundle $`TM`$ admits local symplectic trivializations whose domains are fixed by the diffeomorphisms $`\psi _t`$, we deduce a formula for $`I_\psi `$ in which appear a contribution related to the Maslov indices of the linearized flow $`\psi _t`$ in the trivializations, and a second one in which are involved transition functions of the bundle $`\text{det}(TM)`$. The second contribution is related with the Chern class $`c_1(TM)`$ in the following sense. Using the expression of $`c_1(M)`$ in terms of the transition functions of $`TM`$ determined by the trivializations, $`c_1(M)[\omega ]^{n1},M`$ can be written as a sum $`_j_{R_j}\sigma _j`$, where $`\sigma _j`$ is a $`2n1`$ form (see (3.16)). It turns out that the second contribution is equal to this sum “weighted” by a multiple of the Hamiltonian $`f_t`$ which generates $`\psi `$; more concretely, that contribution is $`n_j𝑑t_{R_j}(f_t\psi _t)\sigma _j`$.
Let $`(M,\omega ,f)`$ be an integrable system such that the points where the integrals of motion are dependent form a set $`P`$ which is union of codimension $`2`$ submanifolds of $`M`$, and such that $`MP`$ is invariant under $`\psi _t`$ and on it there exist action-angle coordinates. Furthermore we assume that there are $`\psi _t`$ invariant Darboux charts which cover $`P`$. Then the expression of $`I_\psi `$ in this atlas reduces to the aforesaid second contribution; that is, $`I_\psi =n_j_{R_j}f\sigma _j`$.
The paper is organized as follows. In Section 2 we recall the construction of the coupling class $`c`$ following . Section 3 is concerned with the proof of the mentioned expression for $`I_\psi `$. First we express $`c_1(M)[\omega ]^{n1},M`$ as the sum $`_j_{R_j}\sigma _j`$ of integrals of $`2n1`$ forms, and next we use this result to prove the formula for the invariant $`I_\psi `$. In Section 4 we check and apply the formulae obtained in Section 3. Using these formulae, we calculate $`I_\psi `$, when $`\psi `$ is the loop in $`\text{Ham}(S^2)`$ generated by the $`1`$-turn rotation of $`S^2`$ around the $`z`$-axis. The result $`I_\psi =0`$ agrees with the fact that $`\pi _1(\text{Ham}(S^2))=_2`$ and $`I`$ is a group homomorphism on $`\text{Ham}(M)`$. We also prove that $`I`$ on $`\pi _1(\text{Ham}(𝕋^{2n}))`$ vanishes identically. When $`n=1`$ this result is consistent with the fact that $`\pi _1(\text{Ham}(𝕋^2))=0`$. Finally we determine the value of $`I`$ on the loops generated by action of $`𝕋^2`$ on a general symplectic Hirzebruch surface (see Theorem 8).
I thank Dusa McDuff for explaining me properties of the Maslov index of the linearized flow, and Eva Miranda for clarifying me some points relative to action-angle variables.
## 2. The coupling class
Let $`(M,\omega )`$ be a compact connected symplectic $`2n`$-manifold. Let $`\psi :S^1=/\text{Ham}(M,\omega )`$ be a loop in the group $`\text{Ham}(M,\omega )`$ at id. By $`X_t`$ is denoted the time-dependent vector field generated by $`\psi _t`$ and $`f_t`$ is the normalized time-dependent Hamiltonian; that is,
$$\frac{d\psi _t}{dt}=X_t\psi _t,\iota _{X_t}\omega =df_t,_Mf_t\omega ^n=0.$$
Given $`ϵ`$, with $`0<ϵ<\pi /2`$, we set
$$D_+^2:=\{pS^2|\mathrm{\hspace{0.17em}0}\theta (p)<\pi /2+ϵ\}$$
$$D_{}^2:=\{pS^2|\pi /2ϵ<\theta (p)\pi \},$$
where $`\theta [0,\pi ]`$ is the polar angle from the z-axis.
Next we construct the Hamiltonian bundle $`E`$ over $`S^2`$ determined by $`\psi `$. First of all we extend $`\psi `$ to a map defined on $`F:=D_+^2D_{}^2`$ by putting $`\psi (\theta ,\varphi )=\psi _t`$, with $`t=\varphi /2\pi `$, with $`\varphi `$ the spherical azimuth angle. We set
$$E=[(D_+^2\times M)(D_{}^2\times M)]/,\text{where}$$
$$(+,p,x)(,p^{},y)\text{iff}\{\begin{array}{cc}p=p^{}F,\hfill & \\ y=\psi _t^1(x),t=\varphi (p)/2\pi .\hfill & \end{array}$$
In this way $`ME\stackrel{\pi }{}S^2`$ is a Hamiltonian bundle over $`S^2`$.
We assume that $`D_\pm ^2`$ are endowed with the orientations induced by the usual one of $`S^2`$ (that is, the orientation of $`S^2`$ as border of the unit ball). We suppose that $`S^1`$ is oriented by $`dt=d\varphi /2\pi `$, that is, $`S^1`$ is oriented as $`D_+`$ . In $`E`$ one considers the orientation induced by the one defined on $`D_+^2\times M`$ by $`d\theta d\varphi \omega ^n`$.
Let $`\alpha `$ be a monotone smooth map $`\alpha :[\pi /2ϵ,\pi ][0,\mathrm{\hspace{0.17em}1}]`$, with $`\alpha (\theta )=1`$ for $`\theta [\pi /2ϵ,\pi /2+ϵ]`$ and $`\alpha (\theta )=0`$ for $`\theta `$ near $`\pi `$. Now we consider the $`2`$-form (see )
(2.1)
$$\tau =\{\begin{array}{cc}\omega ,\text{on}D_+^2\times M\hfill & \\ \omega +d(\alpha (f_t\psi _t))dt,\text{on}D_{}^2\times M.\hfill & \end{array}$$
As $`\alpha `$ vanishes near $`\pi `$, $`\tau `$ is well-defined on $`D_{}^2\times M`$; moreover on $`F\times MD_{}^2\times M`$, $`\tau `$ reduces to $`\omega +d(f_t\psi _t)dt`$. If we denote by $`h`$ the map
$$h:F\times MD_{}^2\times MF\times MD_+^2\times M$$
given by $`h(p,x)=(p,\psi _t(x))`$, with $`t=\varphi (p)/2\pi `$, then taking into account that $`h_{}(\frac{}{t})=\frac{}{t}+X_t\psi _t`$, it follows from $`\iota _{X_t}\omega =df_t`$ that $`h^{}\omega =\omega +d(f_t\psi _t)dt`$. So one has the following Proposition
###### Proposition 1.
$`\tau `$ defines a closed $`2`$-form on $`E`$.
Moreover the cohomology class $`[\tau ]H^2(E,)`$ restricted to each fibre coincides with $`[\omega ]`$. On the other hand
$$_E\tau ^{n+1}=(n+1)_{D_{}^2\times M}(f_t\psi _t)\alpha ^{}(\theta )𝑑\theta dt\omega ^n.$$
From the normalization condition for $`f_t`$ it follows that $`_E\tau ^{n+1}=0`$. Hence $`[\tau ]`$ is the coupling class $`c`$ of the fibration $`E`$ .
## 3. The characteristic number $`I_\psi `$.
Denoting $`TM=\{v_xT_xM|xM\}`$, we put
$$VTE=[(D_+^2\times TM)(D_{}^2\times TM)]/,$$
with
$$(+,p,v_x)(,p^{},v_x^{}^{})\text{iff}p=p^{},x^{}=\psi _t^1(x),v_x^{}^{}=(\psi _t^1)_{}(v_x)$$
where $`t=\varphi (p)/2\pi `$. So $`VTE`$ is a vector bundle over $`E`$; by construction it is the vertical tangent bundle of $`E`$.
Let $`(U;X_1,\mathrm{},X_{2n})`$ be a symplectic trivialization of $`TM`$ on $`UM`$, and $`(V;Y_1,\mathrm{},Y_{2n})`$ be a symplectic trivialization on $`VM`$. We put
(3.1)
$$U_\pm :=\{[\pm ,p,x]|pD_\pm ^2,xU\}$$
and similarly for $`V_\pm `$.Denoting $`x_t:=\psi _t^1(x)`$ one has
$$U_+U_{}=\{[+,p,x]|pF,xU,x_tU\}$$
$$V_+V_{}=\{[+,p,x]|pF,xV,x_tV\}$$
$$V_{}U_{}=\{[,p,x]|pD_{}^2,xVU\}$$
$$U_+V_+=\{[+,p,x]|pD_+^2,xVU\}$$
The corresponding transition functions of $`VTE`$ are
$$g_{U_{}U_+}([+,p,x])=A(t,x)Sp(2n,),\text{with}\psi _t^1\left(X_i(x)\right)=\underset{k}{}A_i^k(t,x)X_k(x_t)$$
$$g_{V_{}V_+}([+,p,x])=B(t,x)Sp(2n,),\text{with}\psi _t^1\left(Y_i(x)\right)=\underset{k}{}B_i^k(t,x)Y_k(x_t)$$
$$g_{U_{}V_{}}([,p,x])=R(x)=g_{U_+V_+}([+,p,x]),\text{with}Y_i(x)=\underset{k}{}R_i^k(x)X_k(x).$$
We denote by $`\rho `$ the usual map $`\rho :Sp(2n,)U(1)`$ which restricts to the determinant map on $`U(n)`$ , then $`l_{ab}:=\rho g_{ab}`$ is a transition function for $`\text{det}(VTE)`$. We also use the following notation, the matrices in $`Sp(2n,)`$ are denoted with capital letters and its images by $`\rho `$ will be denoted by the corresponding small letters; that is,
(3.2)
$$a(t,x):=\rho (A(t,x)),b(t,x):=\rho (B(t,x)),r_{UV}(x):=\rho (R(x)).$$
If $`\psi _t(U)U`$ for all $`t`$, given $`xU`$, the winding number of the map $`tS^1a^1(t,x)U(1)`$ is the integer
(3.3)
$$\frac{i}{2\pi }_0^1a^1(t,x)\frac{a}{t}(t,x)𝑑t.$$
This integer is independent of the point $`xU`$, it will be denoted $`J_U`$. The number $`J_U`$ is the Maslov index in $`U`$ of the linearized flow $`\psi _t`$. Analogously, if $`\psi _t(V)V`$ for all $`t`$ we have the integer
(3.4)
$$J_V=\frac{i}{2\pi }_0^1b^1(t,x)\frac{b}{t}(t,x)𝑑t,$$
$`x`$ being any point of $`V`$; this is the Maslov index in $`V`$ of $`\psi _t`$.
As a previous step to compute $`I_\psi `$ we shall prove the following Lemma, in which the value $`c_1(M)[\omega ]^{n1},[M]`$ is expressed in terms of transition functions of $`\text{det}(TM)`$.
###### Lemma 2.
Let $`\{B_1,\mathrm{},B_m\}`$ be a set of trivializations of $`TM`$, such that its domains cover $`M`$. Then
(3.5)
$$c_1(TM)[\omega ]^{n1},[M]=\frac{i}{2\pi }\underset{i<k}{}_{A_{ik}}d(\mathrm{log}s_{ik})\omega ^{n1},$$
$`s_{ik}`$ being the corresponding transition function of $`\text{det}(TM)`$ and
(3.6)
$$A_{ik}=(B_i_{r<k}B_r)B_k.$$
###### Proof.
$`c_1(M)`$ is represented on $`B_a`$ by the $`2`$-form
$$\frac{i}{2\pi }\underset{c}{}d\left(\phi _cd\mathrm{log}s_{ac}\right),$$
where $`\{\phi _c\}`$ is a partition of unity subordinate to the covering $`\{B_1,\mathrm{},B_m\}`$.
If $`m=2`$
$$c_1(M)[\omega ]^{n1},[M]=\frac{i}{2\pi }_{B_1}d\left(\phi _2d\mathrm{log}s_{12}\right)\omega ^{n1}+\frac{i}{2\pi }_{B_2B_1}d\left(\phi _1d\mathrm{log}s_{21}\right)\omega ^{n1}.$$
By Stokes’ theorem
(3.7)
$$c_1(M)[\omega ]^{n1},[M]=_{B_1}\phi _2L_{12}+_{(B_2B_1)}\phi _1L_{21},$$
where
$$L_{jk}:=(i/2\pi )d\mathrm{log}s_{jk}\omega ^{n1}.$$
Since $`(B_2B_1)B_1=\mathrm{}`$, $`\phi _1`$ vanishes on $`(B_2B_1)`$ and the last integral in (3.7) is zero.
As $`\phi _2`$ is $`1`$ on $`B_1`$, we have
$$c_1(M)[\omega ]^{n1},[M]=_{B_1}L_{12}.$$
In this case $`B_1B_2`$, so $`B_1=A_{12}`$, and the the Lemma is proved when $`m=2`$.
If $`m=3`$
(3.8) $`c_1(M)[\omega ]^{n1},[M]`$ $`={\displaystyle _{B_1}}\left(\phi _2L_{12}+\phi _3L_{13}\right)+{\displaystyle _{(B_2B_1)}}\left(\phi _1L_{21}+\phi _3L_{23}\right)`$
(3.9) $`+{\displaystyle _{(B_3(B_1B_2))}}\left(\phi _1L_{31}+\phi _2L_{32}\right).`$
As $`(B_3(B_1B_2))`$ and the interior of $`B_1B_2`$ are disjoint sets, $`\phi _1`$ and $`\phi _2`$ vanish on $`(B_3(B_1B_2))`$, and the integral in (3.9) is zero. Analogously $`(B_2B_1)`$ and support of $`\phi _1`$ are disjoint so
(3.10)
$$c_1(M)[\omega ]^{n1},[M]=_{B_1}\left(\phi _2L_{12}+\phi _3L_{13}\right)+_{(B_2B_1)}\phi _3L_{23}.$$
On the other hand $`B_1=A+D`$, with $`A:=B_1B_2`$ (oriented as $`B_1`$) and $`D:=(B_1B_2)B_3`$ (see Figure 1). Moreover $`(B_2B_1)=A+C`$ with $`C:=(B_2B_1)B_3`$ (oriented as $`B_2`$).
Since $`C(B_1B_2)=\mathrm{}`$, then $`\phi _3|_C=1`$; thus
(3.11)
$$c_1(M)[\omega ]^{n1},[M]=_{A+D}\left(\phi _2L_{12}+\phi _3L_{13}\right)+_A\phi _3L_{23}+_{A_{23}}L_{23}.$$
The last integral in (3.11) is just the term in (3.5) with $`i=2,k=3`$.
Since $`\phi _j|_D=0`$, for $`j=1,2`$, then $`\phi _3|_D=1`$. As $`A`$ and support of $`\phi _1`$ are disjoint sets, then $`(\phi _2+\phi _3)|_A=1`$. It follows from these facts together with the cocycle condition $`L_{13}+L_{32}=L_{12}`$ that
(3.12)
$$c_1(M)[\omega ]^{n1},[M]=_AL_{12}+_DL_{13}+_{A_{23}}L_{23}.$$
On the other hand $`A_{12}=(B_1B_1)B_2=A`$. Similarly $`A_{13}=D`$. Therefore (3.12) is the formula given in the statement of Lemma when $`m=3`$.
The preceding arguments can be generalized to any $`m`$
(3.13) $`c_1(TM)[\omega ]^{n1},[M]`$ $`={\displaystyle _{B_1}}{\displaystyle \underset{j1}{}}\phi _jL_{1j}+\mathrm{}+{\displaystyle _{(B_{m1}_{r<m1}B_r)}}{\displaystyle \underset{jm1}{}}\phi _jL_{m1,j}`$
(3.14) $`+{\displaystyle _{(B_m_{r<m}B_r)}}{\displaystyle \underset{jm}{}}\phi _jL_{m1,j}.`$
For any $`j=1,\mathrm{},m1`$ support of $`\phi _j`$ and $`(B_m_{r<m}B_r)`$ are disjoint sets. Thus the integral (3.14) is zero (as in the cases $`m=2,3`$). We decompose
$$(B_{m1}_{r<m1}B_r)=E+G,$$
with
$$E:=(B_{m1}_{r<m1}B_r)B_m.$$
Then $`\phi _j|_E=0`$ for all $`jm`$ and $`\phi _m|_E=1`$; thus
(3.15)
$$_{(B_{m1}_{r<m1}B_r)}\underset{jm1}{}\phi _jL_{m1,j}=_G+_{A_{m1,m}}L_{m1,m}.$$
The last integral in (3.15) is the term in (3.5) which corresponds to $`i=m1,k=m`$. An analogous, but more tedious, calculation to the one for the case $`m=3`$ allows to identify in (3.13) the remainder terms of (3.5).
Lemma 2 gives a way for expressing $`c_1(TM)[\omega ]^{n1},[M]`$ as a sum of integrals of $`2n1`$ differential forms on $`2n1`$ chains. The righthand side of (3.5) can be written schematically
(3.16)
$$\underset{j}{}_{R_j}\sigma _j.$$
In next Theorem we use this expression to give an explicit formula for $`I_\psi `$ in terms of transition functions of $`\text{det}(TM)`$ and Maslov indices of $`\psi _t`$.
###### Theorem 3.
If $`\{B_1,\mathrm{},B_m\}`$ is a set of symplectic trivializations for $`TM`$ which covers $`M`$, and such that $`\psi _t(B_j)=B_j`$, for all $`t`$ and all $`j`$, then
(3.17)
$$I_\psi =\underset{i=1}{\overset{m}{}}J_i_{B_i_{j<i}B_j}\omega ^n+\underset{i<k}{}N_{ik},$$
where
$$N_{ik}=n\frac{i}{2\pi }_0^1𝑑t_{A_{ik}}(f_t\psi _t)(d\mathrm{log}r_{ik})\omega ^{n1},$$
$`A_{ik}=(B_i_{r<k}B_r)B_k`$, $`J_i`$ is the Maslov index of $`(\psi _t)_{}`$ in the trivialization $`B_i`$ and $`r_{ik}`$ the corresponding transition function of $`\text{det}(TM)`$.
###### Proof.
Using the notation (3.1) we put
(3.18)
$$O_{\mathrm{𝟐}𝐚\mathrm{𝟏}}:=(B_a)_{},O_{\mathrm{𝟐}𝐚}:=(B_a)_+.$$
Then $`\{O_𝐜|c=1,\mathrm{},2m\}`$ is a covering for $`E`$. We shall denote by $`l_{\mathrm{𝐛𝐜}}`$ the respective transition functions for $`\text{det}(VTE)`$. If we set $`U:=B_1,V:=B_2`$, one has by (3.2)
$$l_{\mathrm{𝟏𝟐}}=a(t,x),l_{\mathrm{𝟏𝟑}}=r_{UV}(x),l_{\mathrm{𝟑𝟒}}=b(t,x).$$
We can determine $`I_\psi =c_1(VTE)c^n,[E]`$ applying the result given in Lemma 2 to the set $`\{O_𝐜\}`$ of trivializations of $`VTE`$. That is,
(3.19)
$$I_\psi =\underset{𝐚<𝐛}{}𝒯_{\mathrm{𝐚𝐛}},\text{where}𝒯_{\mathrm{𝐚𝐛}}=\frac{i}{2\pi }_{A_{\mathrm{𝐚𝐛}}}d\mathrm{log}l_{\mathrm{𝐚𝐛}}\tau ^n.$$
It follows from (3.18) and (2.1) that $`\tau `$ is equal to $`\omega `$ on $`A_{\mathrm{𝐚𝐛}}`$ unless $`𝐚`$ and $`𝐛`$ are both odd; in this case $`\tau =\omega +d(\alpha (f_t\psi _t))dt.`$
We will calculate the summand $`𝒯_{\mathrm{𝟏𝟐}}`$ in (3.19). The set $`A_{\mathrm{𝟏𝟐}}=O_\mathrm{𝟏}O_\mathrm{𝟐}=U_{}U_+`$, and
$$U_{}=\{[+,p,x]|pD_{}^2,xU\}\{[,p,x]|pD_{}^2,xU\}.$$
So
$$A_{\mathrm{𝟏𝟐}}=\{[+,p,x]|pD_{}^2,xU\}.$$
Taking into account (3.3), (3.2) together with the fact that orientations of $`S^1`$ and $`D_{}^2`$ are opposite , we deduce
$$𝒯_{\mathrm{𝟏𝟐}}=\frac{i}{2\pi }_U\left(_{S^1}a^1(t,x)\frac{a(t,x)}{t}𝑑t\right)\omega ^n=J_U_U\omega ^n.$$
Next we consider the term $`𝒯_{\mathrm{𝟑𝟒}}`$. The integration domain is
$$A_{\mathrm{𝟑𝟒}}=\left(V_{}(U_{}U_+)\right)V_+=\{[+,p,x]|pD_{}^2,xVU\}.$$
Hence
$$𝒯_{\mathrm{𝟑𝟒}}=J_V_{VU}\omega ^n.$$
In general,
$`A_{\mathrm{𝟐}𝐣\mathrm{𝟏},\mathrm{𝟐}𝐣}`$ $`=(B_j_{r<j}(B_{r+}B_r))B_{j+}`$
$`=\{[+,p,x]|pD_{}^2,xB_j_{r<j}B_r\}.`$
Hence the term in (3.19) with $`𝐚=\mathrm{𝟐}𝐣\mathrm{𝟏}`$, $`𝐛=\mathrm{𝟐}𝐣`$ gives a contribution to $`I_\psi `$ equal to
(3.20)
$$J_{B_j}_{B_jU_{r<j}B_r}\omega ^n$$
Now we analyze $`𝒯_{\mathrm{𝟏𝟑}}`$.
$$A_{\mathrm{𝟏𝟑}}=\{[,p,x]|pD_{}^2,xUV\}.$$
$`D_{}^2`$ is oriented by the form $`d\theta dt`$, and $`UV`$ is oriented with the orientation of $`U`$. Hence
(3.21) $`𝒯_{\mathrm{𝟏𝟑}}`$ $`={\displaystyle \frac{i}{2\pi }}{\displaystyle _{A_{\mathrm{𝟏𝟑}}}}d\mathrm{log}r_{UV}\left(\omega +d(\alpha (f_t\psi _t))dt\right)^n`$
$`={\displaystyle \frac{ni}{2\pi }}{\displaystyle _{A_{\mathrm{𝟏𝟑}}}}d\mathrm{log}r_{UV}(f_t\psi _t))\alpha ^{}(\theta )d\theta dt\omega ^{n1}`$
$`={\displaystyle \frac{+ni}{2\pi }}{\displaystyle _0^1}𝑑t{\displaystyle _{UV}}(f_t\psi _t)d\mathrm{log}r_{UV}\omega ^{n1}.`$
In general, if $`j<k`$
(3.22)
$$𝒯_{\mathrm{𝟐}𝐣\mathrm{𝟏},\mathrm{𝟐}𝐤\mathrm{𝟏}}=\frac{ni}{2\pi }_0^1𝑑t_{A_{jk}}(f_t\psi _t)d\mathrm{log}r_{jk}\omega ^{n1},$$
where $`A_{jk}`$ is the set defined in Lemma 2.
On the other hand
$$A_{\mathrm{𝟏𝟒}}=\left(U_{}(U_+V_{})\right)V_+=\{[,p,x]|pD_{}^2,xUV\}V_+=\mathrm{}.$$
Thus $`𝒯_{\mathrm{𝟏𝟒}}=0`$. In general, for $`j<k`$ the integration domain $`A_{\mathrm{𝟐}𝐣\mathrm{𝟏},\mathrm{𝟐}𝐤}`$ is of the form
$$(B_j)B_{k+}.$$
In the union $``$ appear the sets $`B_k`$ and $`B_{j+}`$, hence
$$A_{\mathrm{𝟐}𝐣\mathrm{𝟏},\mathrm{𝟐}𝐤}(B_j(B_{j+}B_k))B_{k+},$$
and this set is empty by the same reason that $`A_{\mathrm{𝟏𝟒}}=\mathrm{}`$. Therefore $`𝒯_{\mathrm{𝟐}𝐣\mathrm{𝟏},\mathrm{𝟐}𝐤}=0`$, for any $`j<k`$.
The set $`A_{\mathrm{𝟐𝟑}}`$ is
$$A_{\mathrm{𝟐𝟑}}=(U_+U_{})V_{}=\{[+,p,x]|pF,xUV\}.$$
As $`d\mathrm{log}l_{\mathrm{𝟐𝟑}}\omega ^n`$ does not contain $`d\theta `$, the term $`𝒯_{\mathrm{𝟐𝟑}}`$ vanishes. In general, if $`j<k`$
$`A_{\mathrm{𝟐}𝐣,\mathrm{𝟐}𝐤\mathrm{𝟏}}`$ $`=(B_{j+})B_k(B_{j+}B_j)B_k`$
$`=\{[+,p,x]|pF,xB_jB_k\}.`$
Then $`𝒯_{\mathrm{𝟐}𝐣,\mathrm{𝟐}𝐤\mathrm{𝟏}}`$ vanishes by the same reason that $`𝒯_{\mathrm{𝟐𝟑}}=0.`$
Analogous arguments as the ones explained in the preceding paragraph show that $`𝒯_{\mathrm{𝟐}𝐣,\mathrm{𝟐}𝐤}=0`$, for any $`j<k`$.
So, apart from the terms $`𝒯_{\mathrm{𝐚𝐛}}`$ considered in (3.20) and in (3.22), the remainder summands in (3.19) are zero. The theorem follows from (3.20) and (3.22).
From the definition of product in $`\pi _1(\text{Ham}(M,\omega ))`$ by juxtaposition of paths and under the hypotheses of Theorem 3 is obvious that
$$I:\pi _1(\text{Ham}(M,\omega ))$$
is a group homomorphism. This fact has been proved in for the general case.
###### Corollary 4.
If $`U`$ and $`V`$ are symplectic trivializations of $`TM`$, with $`\psi _t(U)=U`$, $`\psi _t(V)=V`$, for all $`t`$ and $`UV=M`$ and $`_{S^1}(f_t\psi _t)𝑑t`$ is a constant $`k`$ on $`UV`$, then
$$I_\psi =J_U_U\omega ^n+J_V_{VU}\omega ^nnkc_1(TM)[\omega ]^{n1},M.$$
###### Corollary 5.
If $`TM`$ is trivial on $`U:=M\{q\}`$, where $`q`$ is a point of $`M`$ fixed by $`\psi _t`$ for all $`t`$, then
$$I_\psi =J_U_M\omega ^nn\left(_{S^1}f_t(q)𝑑t\right)c_1(TM)[\omega ]^{n1},M.$$
Now we analyze the expression for $`I_\psi `$ given in Theorem 3 in case of integrable systems. Let $`f`$ be the normalized Hamiltonian which generates the loop $`\psi `$. We assume that $`(M,\omega ,f)`$ is completely integrable, with $`f_1=f,f_2,\mathrm{},f_n`$ integrals of motion. We suppose that $`df_1,\mathrm{},df_n`$ are independent at the points of $`MP=:V`$, where $`P`$ is a finite union of $`2n2`$ dimensional submanifolds of $`M`$. We suppose that on $`V`$ are defined action-angle coordinates. We put
$$Q:=\{xP|\text{dim\hspace{0.17em}\hspace{0.17em}Span}(df_1(x),\mathrm{},df_n(x))=n1\}.$$
By $`Q_1,\mathrm{},Q_k`$ are denoted the connected components of $`Q`$, and let $`V_j`$ be a tubular neighborhood of $`Q_j`$ in $`M`$, invariant under $`\psi _t`$ for all $`t`$. We assume that on $`V_j`$ is defined a symplectic trivialization of $`TM`$. Then, for each $`j`$ one can choose a family of tubular neighborhoods $`\{V_{jb}V_j\}_{b=1,2\mathrm{}}`$, such that
$$\underset{b\mathrm{}}{lim}_{V_{jb}}\omega ^n=0.$$
Lemma 2 applied to the covering $`\{V,V_{jb}\}_{j=1,\mathrm{},k}`$ of $`VQ`$ gives
$$c_1(M)[\omega ]^{n1},[M]=\frac{i}{2\pi }\underset{j=1}{\overset{k}{}}_{VV_{jb}}d\mathrm{log}r_{VV_{jb}}+ϵ(b),$$
where $`ϵ(b)`$ goes to $`0`$ as $`b\mathrm{}`$.
Hence
(3.23)
$$c_1(M)[\omega ]^{n1},[M]=\underset{j=1}{\overset{k}{}}z_j,$$
with
(3.24)
$$z_j:=\frac{i}{2\pi }_{VV_{jb}}d\mathrm{log}r_{VV_{jb}}\omega ^{n1}.$$
###### Proposition 6.
Let $`(M,\omega ,f,f_2,\mathrm{}f_n)`$ be an integrable in which the preceding hypotheses hold, then
$$I_\psi =\underset{j=1}{\overset{k}{}}z_j^{},$$
where $`z_j^{}`$ is obtained from the corresponding $`z_j`$ by inserting the factor $`nf`$ in the integrand of (3.24).
###### Proof.
The Maslov index $`J_V=0`$ because of the particular form of the flow equations in action-angle coordinates. On the other hand
$$_{V_{jb}(V\mathrm{})}\omega ^n=0.$$
Thus the Proposition follows from Theorem 3, together with (3.23) and (3.24). ∎
Similar arguments to the ones involved in this Proposition are used in Section 4 for studying the invariant $`I`$ in Hirzebruch surfaces.
## 4. Examples.
The invariant $`I`$ when the manifold is the $`2`$-sphere.
Let $`\psi _t`$ be the rotation in $`^3`$ around $`\stackrel{}{e}_3`$ of angle $`2\pi t`$ with $`t[0,\mathrm{\hspace{0.17em}1}]`$. Then $`\psi _t`$ determines a Hamiltonian symplectomorphism of $`(S^2,\omega _{area})`$. In fact, the isotopy $`\{\psi _t\}`$ is generated by the vector field $`\frac{}{\varphi }`$, and the function $`f`$ on $`S^2`$ defined by $`f(\theta ,\varphi )=2\pi \mathrm{cos}\theta =2\pi z`$ is the corresponding normalized Hamiltonian.
$`TS^2`$ can be trivialized on $`U=D_+^2`$, and on $`V=D_{}^2`$. Moreover $`UV`$ is the parallel $`\theta =\pi /2+ϵ`$. On $`UV`$ the function $`f\psi _t`$ takes the value $`2\pi \mathrm{sin}ϵ`$.
$$_U\omega =2\pi (1k^{}),_{VU}\omega =2\pi (1+k^{}),$$
with $`k^{}:=\mathrm{cos}(\pi /2+ϵ)`$.
Furthermore the north pole $`n`$ and the south pole $`s`$ are fixed points of the isotopy $`\psi _t`$. The rotation $`\psi _t`$ transforms the basis $`\stackrel{}{e}_1,\stackrel{}{e}_2`$ of $`T_nS^2`$ in
$$(\mathrm{cos}2\pi t\stackrel{}{e}_1+\mathrm{sin}2\pi t\stackrel{}{e}_2,\mathrm{sin}2\pi t\stackrel{}{e}_1+\mathrm{cos}2\pi t\stackrel{}{e}_2).$$
So $`J_U`$ is the winding number of the map
$$t[0,\mathrm{\hspace{0.17em}1}]e^{2\pi ti}U(1);$$
That is, $`J_U=+1`$.
Similarly, by considering the oriented basis $`\stackrel{}{e}_2,\stackrel{}{e}_1`$ of $`T_sS^2`$ it turns out that the Maslov index $`J_V`$ of $`\psi _t`$ is $`1`$.
By Corollary 4
$$I_\psi =2\pi (1k^{})2\pi (1+k^{})(2\pi k^{})c_1(TS^2),S^2=0.$$
Corollary 5 can also be applied to determine $`I_\psi `$. One takes $`U:=S^2\{s\}`$. As $`f(s)=2\pi (1)`$, we obtain again
$$I_\psi =+4\pi 2\pi c_1(TS^2),S^2=0.$$
Using formula (LABEL:finalfr) we can determine $`I_\psi `$ again. Now $`V`$ is $`S^2\{n,s\}`$, $`U_1`$ is a small polar cap at $`n`$ and $`U_2`$ the symmetric one at $`s`$. By the symmetry
$$_{U_1V}d\mathrm{log}r_{U_1V}=_{U_2V}d\mathrm{log}r_{U_2V},$$
so $`y_1=y_2`$. As $`f(n)=f(s)`$, then $`I_\psi =0`$.
This result was expected, because $`\pi _1(\text{Ham}(S^2))`$ is isomorphic to $`_2`$ (see ) and $`I`$ is a group homomorphism.
The invariant $`I`$ for Hamiltonian loops in $`𝕋^{2n}`$.
We identify the torus $`𝕋^{2n}`$ with $`^{2n}/^{2n}`$, and we suppose that $`𝕋^{2n}`$ is equipped with the standard symplectic form $`\omega _0`$. If $`\psi _t`$ is a Hamiltonian isotopy of $`𝕋^{2n}`$, it can be written in the form
$$\psi _t(x^1,\mathrm{},x^{2n})=(x^1+\alpha ^1(t,x^i),\mathrm{},x^{2n}+\alpha ^{2n}(t,x^i)),$$
where the function $`\alpha ^j`$, for $`j=1,\mathrm{},2n`$, is periodic of period $`1`$ in each variable: $`t,x^1,\mathrm{},x^{2n}`$. The vector fields $`\{\frac{}{x^i}\}`$ give a symplectic trivialization of the tangent bundle. In this case the right hand side of (3.17) has only one term. The matrix of $`(\psi _t)_{}`$ with respect to $`\{\frac{}{x^i}\}`$ is
(4.1)
$$\left(\delta _i^j+\frac{\alpha ^j}{x^i}\right)Sp(2n,).$$
First, let us assume that each $`\alpha ^j`$ is a separate variables function; that is, $`\alpha ^j(t,x^i)=f^j(t)u^j(x^i)`$. Since $`\alpha ^1`$ takes the same value at symmetric points on opposite faces of the cube $`I^{2n}`$, there is a point $`p_1I^{2n}`$ such
$$\frac{u^1}{x^j}(p_1)=0,$$
for all $`j`$. Hence the first row of the matrix (4.1) at the point $`p_1`$ is $`(1,0,\mathrm{},0)`$; that is, the matrix of $`(\psi _t)_{}(p_1)`$ is independent of $`f^1`$ and thus the Maslov index of $`\{(\psi _t)_{}(p_1)\}_t`$ does not depend on $`f^1`$. From (3.17) it follows that $`I_\psi `$ is independent of $`f^1`$. The independence of $`I_\psi `$ with respect to $`f^i`$ is proved in a similar way. Thus in order to determine $`I_\psi `$ we can assume that $`f^i=0`$ for all $`i`$, but in this case $`I_\psi =0`$ obviously.
If $`\alpha ^j`$ is sum of two separate variables functions
$$\alpha ^j(t,x^i)=f^j(t)u^j(x^i)+g^j(t)v^j(x^i),$$
we take a point $`q_1I^{2n}`$, such that $`\frac{v^1}{x^j}(q_1)=0,`$ for all $`j`$. Then $`I_\psi `$ is independent of $`g^1`$. The above reasoning gives $`I_\psi =0`$ in this case as well.
By the Fourier theory, the original $`C^{\mathrm{}}`$ periodic function $`\alpha ^j`$ can be approximated (in the uniform $`C^k`$-norm) by a sum of separated functions of the form $`f_a(t)u_a(x^i)`$, where $`f_a`$ and $`u_a`$ are $`1`$-periodic. As $`I_\psi `$ depends only on the homotopy class of $`\psi `$, we conclude that $`I_\psi =0`$ for a general Hamiltonian loop.
###### Proposition 7.
The invariant $`I`$ is identically zero on $`\pi _1(\text{Ham}(𝕋^{2n},\omega _0))`$.
This result when $`n=1`$ is consistent with the fact that $`\pi _1(\text{Ham}(𝕋^2))=0`$ (see )
Application to Hirzebruch surfaces.
Given $`3`$ numbers $`k,\tau ,\mu `$, with $`k_{>0}`$, $`\tau ,\mu _{>0}`$ and $`k\mu <\tau `$, the triple $`(k,\tau ,\mu )`$ determine a Hirzebruch surface $`M_{k,\tau ,\mu }`$ . This manifold is the quotient
$$\{z^4:k|z_1|^2+|z_2|^2+|z_4|^2=\tau /\pi ,|z_1|^2+|z_3|^2=\mu /\pi \}/𝕋^2,$$
where the $`𝕋^2`$-action is given by
$$(a,b)(z_1,z_2,z_3,z_4)=(a^kbz_1,az_2,bz_3,az_4),$$
for $`(a,b)𝕋^2`$. The map
$$[z_1,z_2,z_3,z_4]([z_2:z_4],[z_2^kz_3:z_4^kz_3:z_1])$$
allows us to represent $`M_{k,\tau ,\mu }`$ as a submanifold of $`P^1\times P^2`$. On the other hand the usual symplectic structures on $`P^1`$ and $`P^2`$ induce a symplectic form $`\omega `$ on $`M_{k,\tau ,\mu }`$, and the following $`𝕋^2`$-action on $`P^1\times P^2`$
$$(a,b)([u_0:u_1],[x_0:x_1:x_2])=([au_0:u_1],[a^kx_0:x_1:bx_2])$$
gives rise to a toric structure on $`M_{k,\tau ,\mu }`$. In terms of the Delzant construction $`(M_{k,\tau ,\mu },\omega )`$ is associated to the trapezoid in $`(^2)^{}`$ whose not oblique edges are $`\tau ,\mu ,\lambda :=\tau k\mu `$, (see Figure 2). Moreover $`\lambda `$ is the value that the symplectic form $`\omega `$ takes on the exceptional divisor, $`\{[z]M|z_3=0\}`$, of $`M:=M_{k,\tau ,\mu }`$. And $`\omega `$ takes the value $`\mu `$ on the class of the fibre in the fibration $`MP^1`$.
Since $`M`$ is a toric manifold, the $`𝕋^2`$-action define symplectomorphisms of $`M`$. More precisely, let $`\psi _t`$ the diffeomorphism of $`M`$ defined by
(4.2)
$$\psi _t[z_1,z_2,z_3,z_4]=[z_1e^{2\pi it},z_2,z_3,z_4].$$
$`\psi =\{\psi _t:t[0,1]\}`$ is a loop of Hamiltonian symplectomorphisms of $`(M,\omega )`$. Similarly we have
(4.3)
$$\stackrel{~}{\psi }_t[z_1,z_2,z_3,z_4]=[z_1,z_2e^{2\pi it},z_3,z_4],$$
and the corresponding loop $`\stackrel{~}{\psi }`$ in $`\text{Ham}(M,\omega ).`$
Using Theorem 3 we shall calculate the values of $`I_\psi `$ and $`I_{\stackrel{~}{\psi }}`$ in terms of $`\lambda ,`$ $`\tau `$ and $`k`$. The result is stated in Theorem 8 below. The most laborious point in the proof of the following Theorem is to obtain Darboux charts for $`M`$ which give rise to simple transition functions for $`\text{det}(TM)`$.
###### Theorem 8.
Let $`\psi `$ and $`\stackrel{~}{\psi }`$ be the loops of symplectomorphisms of the Hirzebruch surface $`(M_{k,\tau ,\mu },\omega )`$, defined by (4.2) and (4.3) respectively, then
$$I_\psi =\frac{2k\mu ^2}{3}\left(1\frac{\mu }{2\lambda +k\mu }\right),\text{and}I_{\stackrel{~}{\psi }}=\frac{k^2\mu ^2}{3}\left(1\frac{\mu }{2\lambda +k\mu }\right).$$
$`\lambda `$ being $`\tau k\mu `$.
###### Proof.
We will define a Darboux atlas on $`M`$. First we consider the following covering for $`M`$
$$U_1=\{[z]M:z_30z_4\},U_2=\{[z]M:z_10z_4\}$$
$$U_3=\{[z]M:z_10z_2\},U_4=\{[z]M:z_20z_3\}.$$
We set $`z_j=\rho _je^{i\theta _j}`$, with $`\rho _j=|z_j|`$, and on $`U_1`$ introduce the coordinates $`(x_1,y_1,a_1,b_1)`$ by the formulae
$$x_1+iy_1=\rho _1e^{i\phi _1},a_1+ib_1=\rho _2e^{i\phi _2},\phi _1=\theta _1\theta _3k\theta _4,\phi _2=\theta _2\theta _4.$$
Then $`\omega `$ on $`U_1`$ can be written $`\omega =dx_1dy_1+da_1db_1.`$
On $`U_2`$ we consider the Darboux coordinates $`(x_2,y_2,a_2,b_2)`$, with
$$x_2+iy_2=\rho _3e^{i\xi _3},a_2+ib_2=\rho _2e^{i\xi _2},\xi _2=\theta _2\theta _4,\xi _3=\theta _3\theta _1+k\theta _4.$$
On $`U_3`$ we put
$$x_3+iy_3=\rho _3e^{i\chi _3},a_3+ib_3=\rho _4e^{i\chi _4},\chi _3=\theta _3\theta _1+k\theta _2,\chi _4=\theta _4\theta _2,$$
and $`\omega =dx_3dy_3+da_3db_3.`$
Finally, on $`U_4`$ we set
$$x_4+iy_4=\rho _1e^{i\zeta _1},a_4+ib_4=\rho _4e^{i\zeta _4},\zeta _1=\theta _1\theta _3k\theta _2,\zeta _4=\theta _4\theta _2,$$
and $`\omega =dx_4dy_4+da_4db_4.`$
The normalized Hamiltonian function for $`\psi _t`$ is $`f=\pi \rho _1^2\kappa `$, where $`\kappa `$ is a constant determined by the condition $`_Mf\omega ^2=0`$. Straightforward calculations give
$$_M\omega ^2=\mu (2\tau k\mu ),\text{and}_M\pi \rho _1^2\omega ^2=\frac{\mu ^2}{3}(3\tau 2k\mu ).$$
So
(4.4)
$$\kappa =\frac{\mu }{3}\left(\frac{3\lambda +k\mu }{2\lambda +k\mu }\right).$$
It is not easy to determine the transition function of $`\text{det}(TM)`$ that corresponds to the coordinate transformation $`(x_i,y_i,a_i,b_i)(x_j,y_j,a_j,b_j)`$; that is why we will introduce polar coordinate on subsets of the domains $`U_j`$.
Given $`0<ϵ<<1`$, for $`j=1,2,3,4`$ we put
$$B_j=\{[z]U_j:|z_j|<2ϵ\}\text{and}B_0=\{[z]M:|z_j|>ϵ\text{for all}j\}.$$
On $`B_0`$ are well-defined the coordinates $`(\frac{\rho _1^2}{2},\phi _1,\frac{\rho _2^2}{2},\phi _2)`$, and in this coordinates
$$\omega =d\left(\frac{\rho _1^2}{2}\right)d\phi _1+d\left(\frac{\rho _2^2}{2}\right)d\phi _2.$$
On $`B_j`$ ($`j=1,2,3,4`$) we consider the Darboux coordinates $`(x_j,y_j,a_j,b_j)`$ defined above. Then $`B_0,B_1,B_2,B_3,B_4`$ is a Darboux atlas for $`M`$. We assume that $`M`$ is endowed with the orientation given by $`\omega ^2`$. This orientation agrees on $`B_0`$ with the one defined by $`d\rho _1^2d\phi _1d\rho _2^2d\phi _2`$.
It is evident that $`\psi _t(B_i)=B_i`$, for $`i=0,1,2,3,4`$. Since $`\psi _t`$ on $`B_0`$ is simply the translation $`\phi _1\phi _1+2\pi t`$ of the variable $`\phi _1`$, the Maslov index $`J_0`$ of $`\psi `$ in the trivialization defined on $`B_0`$ vanishes.
As $`B_j`$ (for $`j=1,2,3,4`$) has ”infinitesimal size” and $`J_0=0`$, the expression for $`I_\psi `$ of Theorem 3 can be written
(4.5)
$$I_\psi =\underset{i<k}{}N_{ik}+O(ϵ)$$
Since $`I_\psi `$ is obviously independent of the coordinates, it follows from (4.5) that $`N_{ik}`$ is independent, up to order $`ϵ`$, of the chosen Darboux coordinates in $`B_j`$, for $`j=1,2,3,4`$. Moreover $`N_{ik}`$ with $`0i<k`$ is also of order $`ϵ`$.
On the other hand, if we substitute $`B_j`$ by
$$B_j^{}=\{[z]B_j:|z_r|>ϵ,rj\}$$
in the definition of $`N_{ik}`$ (see Theorem 3) the new $`N_{ik}`$ differs from the old one in a quantity of order $`ϵ`$. As on $`B_1^{}`$ the variable $`\rho _20`$, we can consider the Darboux coordinates
$$(x_1,y_1,\frac{\rho _2^2}{2},\phi _2)$$
on $`B_1^{}`$. Since $`\rho _30`$ on $`B_2^{}`$ we take the coordinates $`(a_2,b_2,\frac{\rho _3^2}{2},\xi _3)`$ on $`B_2^{}`$. Similarly we will adopt the following coordinates: $`(x_3,y_3,\frac{\rho _4^2}{2},\chi _4)`$ on $`B_3^{}`$ and $`(\frac{\rho _1^2}{2},\zeta _1,a_4,b_4)`$ on $`B_4^{}`$.
Taking into account the preceding arguments
(4.6)
$$I_\psi =\underset{j=1}{\overset{4}{}}N_{0j}^{}+O(ϵ),$$
where
(4.7)
$$N_{0j}^{}=\frac{i}{\pi }_{A_{0j}^{}}fd\mathrm{log}r_{0j}\omega $$
and
$$A_{0j}^{}=\{[z]M:|z_r|>ϵ,\text{for all}rj\text{and}|z_j|=ϵ\}.$$
The submanifold $`A_{0j}^{}`$ is oriented as a subset of $`B_0`$; that is, with the orientation induced by the one of $`B_0`$.
Next we determine the value of $`N_{01}^{}`$. To know the transition function $`r_{01}`$ one needs the Jacobian matrix $`R`$ of the transformation
$$(x_1,y_1,\frac{\rho _2^2}{2},\phi _2)(\frac{\rho _1^2}{2},\phi _1,\frac{\rho _2^2}{2},\phi _2)$$
in the points of $`A_{01}^{}`$; with $`\rho _1^2=x_1^2+y_1^2`$, $`\phi _1=\mathrm{tan}^1(y_1/x_1)`$. The non trivial block of $`R`$ is the diagonal one
$$\left(\begin{array}{cc}x_1& y_1\\ r& s\end{array}\right),$$
with $`r=y_1(x_1^2+y_1^2)^1`$ and $`s=x_1(x_1^2+y_1^2)^1`$. The non-real eigenvalues of $`R`$ are
$$\lambda _\pm =\frac{x_1+s}{2}\pm \frac{i\sqrt{4(s+x_1)^2}}{2}.$$
On $`A_{01}^{}`$ these non-real eigenvalues occur when $`(s+x_1)^2<2`$, that is, if $`|\mathrm{cos}\phi _1|<2ϵ(ϵ^2+1)^1=:\delta `$. If $`y_1>0`$ then $`\lambda _{}`$ of the first kind (see ) and $`\lambda _+`$ is of the first kind, if $`y_1<0`$.
Hence, on $`A_{01}^{}`$,
$$\rho (R)=\{\begin{array}{cc}\lambda _+|\lambda _+|^1=x+iy,\hfill & \text{if }|\mathrm{cos}\phi _1|<\delta \text{ and }y_1<0\text{;}\hfill \\ \lambda _{}|\lambda _{}|^1=xiy,\hfill & \text{if }|\mathrm{cos}\phi _1|<\delta \text{ and }y_1>0\text{;}\hfill \\ \pm 1,\hfill & \text{otherwise.}\hfill \end{array}$$
where $`x=\delta ^1\mathrm{cos}\phi _1`$, and $`y=\sqrt{1x^2}`$.
If we put $`\rho (R)=e^{i\gamma }`$, then $`\mathrm{cos}\gamma =\delta ^1\mathrm{cos}\phi _1`$ (when $`|\mathrm{cos}\phi _1|<\delta `$), and
$$\mathrm{sin}\gamma =\{\begin{array}{cc}\sqrt{1\mathrm{cos}^2\gamma },\hfill & \text{if }\mathrm{sin}\phi _1>0\text{;}\hfill \\ \sqrt{1\mathrm{cos}^2\gamma },\hfill & \text{if }\mathrm{sin}\phi _1<0\text{.}\hfill \end{array}$$
So when $`\phi _1`$ runs anticlockwise from $`0`$ to $`2\pi `$, $`\gamma `$ goes round clockwise the circumference; that is, $`\gamma =h(\phi _1)`$, where $`h`$ is a function such that
(4.8)
$$h(0)=2\pi ,\text{and}h(2\pi )=0.$$
As $`r_{01}=\rho (R)`$, then $`dlogr_{01}=idh`$.
On $`A_{10}^{}`$ the form $`\omega `$ reduces to $`(1/2)d\rho _2^2d\phi _2`$. From (4.7) one deduces
(4.9)
$$N_{01}^{}=\frac{i}{2\pi }_{A_{01}^{}}if𝑑hd\rho _2^2d\phi _2.$$
On the other hand according to the convention about orientations, $`\{[z]:|z_1|=ϵ\}`$ as subset of $`B_0`$ is oriented by $`d\phi _1d\rho _2^2d\phi _2.`$ And on $`A_{01}^{}`$ the Hamiltonian function $`f=\kappa +O(ϵ)`$. Then it follows from (4.9) together with (4.8)
(4.10)
$$N_{01}^{}=2\tau \kappa +O(ϵ).$$
The contributions $`N_{02}^{},N_{03}^{},N_{04}^{}`$ to $`I_\psi `$ can be calculated in a similar way. One obtains the following results up to addends of order $`ϵ`$:
(4.11)
$$N_{02}^{}=2\mu \kappa \mu ^2,N_{03}^{}=2\lambda (\kappa \mu ),N_{04}^{}=\mu (2\kappa \mu ).$$
As $`I_\psi `$ is independent of $`ϵ`$, it follows from (4.6), (4.10), (4.11) and (4.4)
$$I_\psi =\frac{2k\mu ^2}{3}\left(1\frac{\mu }{2\lambda +k\mu }\right).$$
Next we consider the loop $`\stackrel{~}{\psi }`$; the corresponding normalized Hamiltonian function is $`\stackrel{~}{f}=\pi \rho _2^2\stackrel{~}{\kappa }`$, where
(4.12)
$$\stackrel{~}{\kappa }=\frac{3\lambda ^2+3k\lambda \mu +k^2\mu ^2}{3(2\lambda +k\mu )}.$$
As in the preceding case
(4.13)
$$I_{\stackrel{~}{\psi }}=\underset{j=1}{\overset{4}{}}\stackrel{~}{N}_{0j}^{}+O(ϵ),$$
where
$$\stackrel{~}{N}_{0j}^{}=\frac{i}{\pi }_{A_{0j}^{}}\stackrel{~}{f}d\mathrm{log}r_{0j}\omega .$$
The expression for $`\stackrel{~}{N}_{01}^{}`$ can be obtained from (4.9) substituting $`f`$ for $`\stackrel{~}{f}`$; so
(4.14)
$$\stackrel{~}{N}_{01}^{}=\tau (2\stackrel{~}{\kappa }\tau )+O(ϵ).$$
Analogous calculations give the following values for the $`\stackrel{~}{N}_{0j}^{}`$’s, up to summands of order $`ϵ`$
(4.15)
$$\stackrel{~}{N}_{02}^{}=2\mu \stackrel{~}{\kappa },\stackrel{~}{N}_{03}^{}=\lambda (2\stackrel{~}{\kappa }\lambda ),\stackrel{~}{N}_{04}^{}=\mu (2\stackrel{~}{\kappa }k\mu 2\lambda ).$$
From (4.12), (4.13), (4.14) and (4.15) it follows the value for $`I_{\stackrel{~}{\psi }}`$ given in the statement of Theorem.
Remark. In is proved that $`\pi _1(\text{Ham}(M))=`$ when $`k=1`$, therefore the quotient of $`I_\psi `$ by $`I_\psi ^{}`$, for arbitrary Hamiltonian loops of symplectomorphisms, is a rational number. For the particular loops considered in Theorem 8 the quotient $`I_{\stackrel{~}{\psi }}/I_\psi `$ equals $`k/2`$, so Theorem 8 is consistent with the result of Abreu and McDuff. |
warning/0506/math0506219.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Let $`𝕂`$ denote a field, and let $`V`$ denote a vector space over $`𝕂`$ with finite positive dimension. We consider an ordered pair of linear transformations $`A:VV`$ and $`A^{}:VV`$ that satisfy the following two conditions:
* There exists a basis for $`V`$ with respect to which the matrix representing $`A`$ is irreducible tridiagonal and the matrix representing $`A^{}`$ is diagonal.
* There exists a basis for $`V`$ with respect to which the matrix representing $`A^{}`$ is irreducible tridiagonal and the matrix representing $`A`$ is diagonal.
Such a pair is called a Leonard pair on $`V`$. This notion was introduced by the second author .
Throughout this paper, we fix the following notation. Let $`A`$, $`A^{}`$ denote a Leonard pair on $`V`$. We set $`d=dimV1`$. Let $`v_0^{}`$, $`v_1^{}`$, …, $`v_d^{}`$ denote a basis for $`V`$ that satisfies the condition (i), and let $`v_0`$, $`v_1`$, …, $`v_d`$ denote a basis for $`V`$ that satisfies (ii). For $`0id`$, let $`a_i`$ denote the coefficient of $`v_i^{}`$, when we write $`Av_i^{}`$ as a linear combination of $`v_0^{}`$, $`v_1^{}`$, …, $`v_d^{}`$, and let $`a_i^{}`$ denote the coefficient of $`v_i`$, when we write $`A^{}v_i`$ as a linear combination of $`v_0`$, $`v_1`$, …, $`v_d`$.
In this paper we prove the following results.
###### Theorem 1.1
The following are equivalent.
* $`a_0=a_d`$,
* $`a_0^{}=a_d^{}`$.
###### Theorem 1.2
For $`d1`$ the following are equivalent.
* $`a_0=a_d`$ and $`a_1=a_{d1}`$,
* $`a_0^{}=a_d^{}`$ and $`a_1^{}=a_{d1}^{}`$,
* $`a_i=a_{di}`$ and $`a_i^{}=a_{di}^{}`$ for $`0id`$.
We say that $`A`$, $`A^{}`$ is balanced whenever $`a_i=a_{di}`$ and $`a_i^{}=a_{di}^{}`$ for $`0id`$.
###### Remark 1.3
Theorems 1.1 and 1.2 give a proof of a conjecture by the second author \[5, Section 36\].
###### Remark 1.4
Pascasio \[1, Corollary 4.3\] proved Theorem 1.1 for the Leonard pairs that come from a $`Q`$-polynomial distance-regular graph.
For $`0id`$, let $`\theta _i`$ (respectively $`\theta _i^{}`$) denote the eigenvalue for $`A`$ associated with the eigenvector $`v_i`$ (respectively $`v_i^{}`$). Let $`\phi _1`$, $`\phi _2`$, …, $`\phi _d`$ (respectively $`\varphi _1`$, $`\varphi _2`$, …, $`\varphi _d`$) denote the first split sequence (respectively the second split sequence) with respect to the ordering $`(\theta _0`$, $`\theta _1`$, …, $`\theta _d`$; $`\theta _0^{}`$, $`\theta _1^{}`$, …, $`\theta _d^{})`$. The definition of the split sequences will be given in Section 2.
A Leonard pair is said to be bipartite whenever $`a_i=0`$ for $`0id`$. We consider a slightly more general situation.
###### Theorem 1.5
The following are equivalent.
* $`a_i`$ is independent of $`i`$ for $`0id`$.
* $`\theta _i+\theta _{di}`$ is independent of $`i`$ for $`0id`$, and $`\phi _i=\varphi _i`$ for $`1id`$.
Suppose (i), (ii) hold. Then the common value of $`\theta _i+\theta _{di}`$ is twice the common value of the $`a_i`$.
We say the Leonard pair $`A`$, $`A^{}`$ is essentially bipartite whenever the equivalent conditions (i), (ii) hold in Theorem 1.5. Observe that if $`A`$, $`A^{}`$ is essentially bipartite, then the Leonard pair $`A\xi I`$, $`A^{}`$ is bipartite, where $`\xi `$ denotes the common value of $`a_0`$, $`a_1`$, …, $`a_d`$.
A Leonard pair is said to be dual bipartite whenever $`a_i^{}=0`$ for $`0id`$. We consider a slightly more general situation.
###### Theorem 1.6
The following are equivalent.
* $`a_i^{}`$ is independent of $`i`$ for $`0id`$.
* $`\theta _i^{}+\theta _{di}^{}`$ is independent of $`i`$ for $`0id`$, and $`\phi _i=\varphi _{di+1}`$ for $`1id`$.
Suppose (i), (ii) hold. Then the common value of $`\theta _i^{}+\theta _{di}^{}`$ is twice the common value of the $`a_i^{}`$.
We say the Leonard pair $`A`$, $`A^{}`$ is essentially dual bipartite whenever the equivalent conditions (i), (ii) hold in Theorem 1.6. Observe that if $`A`$, $`A^{}`$ is essentially dual bipartite, then the Leonard pair $`A`$, $`A^{}\xi ^{}I`$ is dual bipartite, where $`\xi ^{}`$ denotes the common value of $`a_0^{}`$, $`a_1^{}`$, …, $`a_d^{}`$.
###### Theorem 1.7
Let $`A`$, $`A^{}`$ denote a Leonard pair.
* If $`A`$, $`A^{}`$ is essentially bipartite, then $`A`$, $`A^{}`$ is balanced.
* If $`A`$, $`A^{}`$ is essentially dual bipartite, then $`A`$, $`A^{}`$ is balanced.
* Assume $`d2`$. If $`A`$, $`A^{}`$ is balanced, then $`A`$, $`A^{}`$ is essentially bipartite or essentially dual bipartite.
###### Remark 1.8
For $`d=2`$, part (iii) of Theorem 1.7 is false. A counter example is given in Example 5.3.
###### Remark 1.9
In our proof of Theorems 1.2, 1.51.7 we use a case-analysis based on the classification of Leonard pairs by the second author .
Our paper is organized as follows. In Section 2 we give some background information. In Section 3 we give the proof of Theorem 1.1. In Section 4 we describe the cases that we will use in our proof of Theorems 1.2, 1.51.7. In Sections 5–10 we give the proofs of these theorems.
## 2 Some background information
In this section we summarize some results that we will use in our proof.
###### Lemma 2.1
\[2, Lemma 1.3\] The eigenvalues $`\theta _0`$, $`\theta _1`$, …, $`\theta _d`$ of $`A`$ are distinct and contained in $`𝕂`$. Moreover, the eigenvalues $`\theta _0^{}`$, $`\theta _1^{}`$, …, $`\theta _d^{}`$ of $`A^{}`$ are distinct and contained in $`𝕂`$.
###### Lemma 2.2
\[2, Lemma 9.5\] For $`d1`$ and for $`0id`$,
$$\frac{\theta _i\theta _{di}}{\theta _0\theta _d}=\frac{\theta _i^{}\theta _{di}^{}}{\theta _0^{}\theta _d^{}}.$$
(1)
###### Theorem 2.3
\[2, Theorem 3.2\] There exists a basis for $`V`$ with respect to which the matrices representing $`A`$, $`A^{}`$ take the following form for some scalars $`\phi _1`$, $`\phi _2`$, …, $`\phi _d`$ in $`𝕂`$:
$$A:\left(\begin{array}{cccccc}\theta _0& & & & & \text{0}\\ 1& \theta _1\\ & 1& \theta _2\\ & & & \\ & & & & \\ \text{0}& & & & 1& \theta _d\end{array}\right),A^{}:\left(\begin{array}{cccccc}\theta _0^{}& \phi _1& & & & \text{0}\\ & \theta _1^{}& \phi _2\\ & & \theta _2^{}& \\ & & & & \\ & & & & & \phi _d\\ \text{0}& & & & & \theta _d^{}\end{array}\right).$$
The sequence $`\phi _1`$, $`\phi _2`$, …, $`\phi _d`$ is uniquely determined by the ordering $`(\theta _0`$, $`\theta _1`$, …, $`\theta _d`$; $`\theta _0^{}`$, $`\theta _1^{}`$, …, $`\theta _d^{})`$. Moreover $`\phi _i0`$ for $`1id`$.
The sequence $`\phi _1`$, $`\phi _2`$, …, $`\phi _d`$ is called the first split sequence with respect to the ordering $`(\theta _0`$, $`\theta _1`$, …, $`\theta _d`$; $`\theta _0^{}`$, $`\theta _1^{}`$, …, $`\theta _d^{})`$. Let $`\varphi _1`$, $`\varphi _2`$, …, $`\varphi _d`$ denote the first split sequence with respect to the ordering $`(\theta _d`$, $`\theta _{d1}`$, …, $`\theta _0`$; $`\theta _0^{}`$, $`\theta _1^{}`$, …, $`\theta _d^{})`$. We call $`\varphi _1`$, $`\varphi _2`$, …, $`\varphi _d`$ the second split sequence with respect to the ordering $`(\theta _0`$, $`\theta _1`$, …, $`\theta _d`$; $`\theta _0^{}`$, $`\theta _1^{}`$, …, $`\theta _d^{})`$. The sequence
$$(\theta _0,\theta _1,\mathrm{},\theta _d;\theta _0^{},\theta _1^{},\mathrm{},\theta _d^{};\phi _1,\phi _2,\mathrm{},\phi _d;\varphi _1,\varphi _2,\mathrm{},\varphi _d)$$
is called a parameter array of the Leonard pair.
In the classification of Leonard pairs, the following theorem plays a key role.
###### Theorem 2.4
\[2, Theorem 1.9\] Let
$$(\theta _0,\theta _1,\mathrm{},\theta _d;\theta _0^{},\theta _1^{},\mathrm{},\theta _d^{};\phi _1,\phi _2,\mathrm{},\phi _d;\varphi _1,\varphi _2,\mathrm{},\varphi _d)$$
(2)
denote a sequence of scalars taken from $`𝕂`$. Then there exists a Leonard pair with parameter array (2) if and only if (i)–(v) hold below.
* $`\phi _i0`$, $`\varphi _i0`$ $`(1id)`$.
* $`\theta _i\theta _j`$, $`\theta _i^{}\theta _j^{}`$ if $`ij`$ $`(0i,jd`$).
* For $`1id`$,
$$\phi _i=\varphi _1\underset{h=0}{\overset{i1}{}}\frac{\theta _h\theta _{dh}}{\theta _0\theta _d}+(\theta _i^{}\theta _0^{})(\theta _{i1}\theta _d).$$
* For $`1id`$,
$$\varphi _i=\phi _1\underset{h=0}{\overset{i1}{}}\frac{\theta _h\theta _{dh}}{\theta _0\theta _d}+(\theta _i^{}\theta _0^{})(\theta _{di+1}\theta _0).$$
* The expressions
$$\frac{\theta _{i2}\theta _{i+1}}{\theta _{i1}\theta _i},\frac{\theta _{i2}^{}\theta _{i+1}^{}}{\theta _{i1}^{}\theta _i^{}}$$
(3)
are equal and independent of $`i`$ for $`2id1`$.
The scalars $`a_i`$, $`a_i^{}`$ can be expressed in terms of the parameter array as follows.
###### Lemma 2.5
\[3, Lemma 10.3\] For $`0id`$,
$$a_i=\theta _i+\frac{\phi _i}{\theta _i^{}\theta _{i1}^{}}\frac{\phi _{i+1}}{\theta _{i+1}^{}\theta _i^{}},a_i^{}=\theta _i^{}+\frac{\phi _i}{\theta _i\theta _{i1}}\frac{\phi _{i+1}}{\theta _{i+1}\theta _i},$$
(4)
$$a_i=\theta _{di}+\frac{\varphi _i}{\theta _i^{}\theta _{i1}^{}}\frac{\varphi _{i+1}}{\theta _{i+1}^{}\theta _i^{}},a_i^{}=\theta _{di}^{}+\frac{\varphi _{di+1}}{\theta _i\theta _{i1}}\frac{\varphi _{di}}{\theta _{i+1}\theta _i},$$
(5)
where we set $`\phi _0=0`$, $`\phi _{d+1}=0`$, $`\varphi _0=0`$, $`\varphi _{d+1}=0`$, and let $`\theta _1`$, $`\theta _{d+1}`$, $`\theta _1^{}`$, $`\theta _{d+1}^{}`$ denote indeterminates.
## 3 Proof of Theorem 1.1
In this section we prove Theorem 1.1.
###### Lemma 3.1
For $`d1`$,
$`a_0`$ $`=`$ $`\theta _0+{\displaystyle \frac{\phi _1}{\theta _0^{}\theta _1^{}}},`$ (6)
$`a_d`$ $`=`$ $`{\displaystyle \frac{\theta _1(\theta _0^{}\theta _d^{})\theta _0(\theta _0^{}\theta _{d1}^{})}{\theta _{d1}^{}\theta _d^{}}}{\displaystyle \frac{\phi _1}{\theta _{d1}^{}\theta _d^{}}},`$ (7)
$`a_0^{}`$ $`=`$ $`\theta _0^{}+{\displaystyle \frac{\phi _1}{\theta _0\theta _1}},`$ (8)
$`a_d^{}`$ $`=`$ $`{\displaystyle \frac{\theta _1^{}(\theta _0\theta _d)\theta _0^{}(\theta _0\theta _{d1})}{\theta _{d1}\theta _d}}{\displaystyle \frac{\phi _1}{\theta _{d1}\theta _d}}.`$ (9)
Proof. The equations (6) and (8) follow from (4). From Theorem 2.4 (iii), (iv),
$$\phi _d=\phi _1+(\theta _1^{}\theta _0^{})(\theta _d\theta _0)+(\theta _d^{}\theta _0^{})(\theta _{d1}\theta _d).$$
(10)
From (1) at $`i=1`$,
$$\theta _d=\theta _0\frac{(\theta _0^{}\theta _d^{})(\theta _1\theta _{d1})}{\theta _1^{}\theta _{d1}^{}}.$$
(11)
Evaluating the equation on the left in (4) using (10) and (11) we find (7). The proof of (9) is similar. $`\mathrm{}`$
###### Lemma 3.2
For $`d1`$,
$`a_0a_d`$ $`=`$ $`{\displaystyle \frac{(\theta _0\theta _1)(\theta _0^{}\theta _d^{})}{\theta _{d1}^{}\theta _d^{}}}+{\displaystyle \frac{\phi _1}{\theta _0^{}\theta _1^{}}}+{\displaystyle \frac{\phi _1}{\theta _{d1}^{}\theta _d^{}}},`$ (12)
$`a_0^{}a_d^{}`$ $`=`$ $`{\displaystyle \frac{(\theta _0^{}\theta _1^{})(\theta _0\theta _d)}{\theta _{d1}\theta _d}}+{\displaystyle \frac{\phi _1}{\theta _0\theta _1}}+{\displaystyle \frac{\phi _1}{\theta _{d1}\theta _d}}.`$ (13)
Proof. Follows from Lemma 3.1. $`\mathrm{}`$
###### Lemma 3.3
For $`d1`$,
$$\frac{(a_0a_d)(\theta _0^{}\theta _1^{})(\theta _{d1}^{}\theta _d^{})}{\theta _0^{}\theta _d^{}}=\frac{(a_0^{}a_d^{})(\theta _0\theta _1)(\theta _{d1}\theta _d)}{\theta _0\theta _d}.$$
(14)
Proof. Using (12) and (13), the left side of (14) becomes
$$\phi _1+(\theta _0\theta _1)(\theta _0^{}\theta _1^{})\frac{\phi _1(\theta _1^{}\theta _{d1}^{})}{\theta _0^{}\theta _d^{}},$$
and the right side of (14) becomes
$$\phi _1+(\theta _0^{}\theta _1^{})(\theta _0\theta _1)\frac{\phi _1(\theta _1\theta _{d1})}{\theta _0\theta _d}.$$
These expressions coincide by (1). $`\mathrm{}`$
Proof of Theorem 1.1 Assume $`d1`$; otherwise the result is vacuously true. Now the result follows from Lemma 3.3. $`\mathrm{}`$
## 4 Description of the cases
Let $`\overline{𝕂}`$ denote the algebraic closure of $`𝕂`$. In our proof of Theorems 1.2, 1.5, 1.6 and 1.7, we break the argument into the following cases.
* Case 0: $`d2`$.
For $`d3`$ let $`q`$ denote a nonzero scalar in $`\overline{𝕂}`$ such that $`q+q^1+1`$ is equal to the common value of (3).
* Case I: $`d3`$, $`q1`$, $`q1`$.
* Case II: $`d3`$, $`q=1`$, $`\text{Char}(𝕂)2`$.
* Case III: $`d3`$, $`q=1`$, $`\text{Char}(𝕂)2`$, $`d`$ even.
* Case IV: $`d3`$, $`q=1`$, $`\text{Char}(𝕂)2`$, $`d`$ odd.
* Case V: $`d3`$, $`q=1`$, $`\text{Char}(𝕂)=2`$.
###### Definition 4.1
For $`d1`$ we let $`H`$ denote the value of (14);
$$H=\frac{(a_0a_d)(\theta _0^{}\theta _1^{})(\theta _{d1}^{}\theta _d^{})}{\theta _0^{}\theta _d^{}}=\frac{(a_0^{}a_d^{})(\theta _0\theta _1)(\theta _{d1}\theta _d)}{\theta _0\theta _d}.$$
We note that $`H=0`$ if and only if $`a_0=a_d`$ if and only if $`a_0^{}=a_d^{}`$.
## 5 Case 0: $`d2`$
In this section we prove Theorems 1.2, 1.51.7 for $`d2`$. We first note that Theorem 1.2 follows from Theorem 1.1 for these values of $`d`$. We consider Theorems 1.51.7.
First assume $`d=0`$. Then Theorems 1.5, 1.6 and 1.7 are vacuously true.
Next assume $`d=1`$. From (6) and the equation on the left in (5) for $`i=1`$,
$$a_0a_1=\frac{\phi _1+\varphi _1}{\theta _0^{}\theta _1^{}}.$$
Thus $`a_0=a_1`$ if and only if $`\phi _1+\varphi _1=0`$. From (6) and (7), we find $`a_0+a_1=\theta _0+\theta _1`$. So that $`2a_0=\theta _0+\theta _1`$ when $`a_0=a_1`$. These imply Theorem 1.5. The proof of Theorem 1.6 is similar. Theorem 1.7 follows from Theorem 1.1.
For the rest of this section, we assume $`d=2`$.
###### Lemma 5.1
The following hold.
$`\phi _1`$ $`=`$ $`H(\theta _0\theta _1)(\theta _0^{}\theta _1^{}),`$ (15)
$`\phi _2`$ $`=`$ $`H(\theta _1\theta _2)(\theta _1^{}\theta _2^{}),`$ (16)
$`\varphi _1`$ $`=`$ $`H+(\theta _1\theta _2)(\theta _0^{}\theta _1^{}),`$ (17)
$`\varphi _2`$ $`=`$ $`H+(\theta _0\theta _1)(\theta _1^{}\theta _2^{}).`$ (18)
Proof. Setting $`d=2`$ in (12) we find (15). The other equations follow from (15) using Theorem 2.4 (iii), (iv). $`\mathrm{}`$
###### Lemma 5.2
Suppose $`H=0`$. Then
$$a_1a_0=\theta _02\theta _1+\theta _2.$$
Proof. Obtained by evaluating the equation on the left in (4) for $`i=0`$, $`1`$ using (15) and (16). $`\mathrm{}`$
Proof of Theorem 1.5 (i)$``$(ii): By assumption $`a_0=a_2`$ so $`H=0`$. Using $`H=0`$ and $`a_0=a_1`$ we find $`\theta _0+\theta _2=2\theta _1`$ by Lemma 5.2. Evaluating the data in Lemma 5.1 using these equations we find $`\phi _1=\varphi _1`$ and $`\phi _2=\varphi _2`$.
(ii)$``$(i): Observe Char$`(𝕂)2`$; otherwise the equation $`\theta _0+\theta _2=2\theta _1`$ becomes $`\theta _0=\theta _2`$ for a contradiction. Comparing (15), (17) we find $`2H=0`$ so $`H=0`$. By this and Definition 4.1 we find $`a_0=a_2`$. Evaluating Lemma 5.2 using $`H=0`$ and $`\theta _0+\theta _2=2\theta _1`$ we find $`a_0=a_1`$. Now $`a_0=a_1=a_2`$ as desired.
Suppose (i), (ii) hold. Evaluating (6) using (15) we find $`a_0=\theta _1`$, so that the common value of $`\theta _i+\theta _{di}`$ is $`2a_0`$. $`\mathrm{}`$
Proof of Theorem 1.6 Similar to the proof of Theorem 1.5. $`\mathrm{}`$
Proof of Theorem 1.7 Follows from Theorem 1.2. $`\mathrm{}`$
We finish this section by giving an example that shows Theorem 1.7 (iii) is false for $`d=2`$.
###### Example 5.3
Let $`\theta _0`$, $`\theta _1`$, $`\theta _2`$, $`\theta _0^{}`$, $`\theta _1^{}`$, $`\theta _2^{}`$ denote scalars in $`𝕂`$ such that $`\theta _i\theta _j`$, $`\theta _i^{}\theta _j^{}`$ if $`ij`$ $`(0i,j2)`$. We define scalars
$`\phi _1`$ $`=`$ $`(\theta _0\theta _1)(\theta _0^{}\theta _1^{}),`$
$`\phi _2`$ $`=`$ $`(\theta _1\theta _2)(\theta _1^{}\theta _2^{}),`$
$`\varphi _1`$ $`=`$ $`(\theta _1\theta _2)(\theta _0^{}\theta _1^{}),`$
$`\varphi _2`$ $`=`$ $`(\theta _0\theta _1)(\theta _1^{}\theta _2^{}).`$
Observe that the sequence
$$(\theta _0,\theta _1,\theta _2;\theta _0^{},\theta _1^{},\theta _2^{};\phi _1,\phi _2;\varphi _1,\varphi _2)$$
(19)
satisfies the conditions (i)–(v) in Theorem 2.4, so that there exists a Leonard pair having the parameter array (19). Using (4), we get
$$a_0=\theta _1,a_1=\theta _0\theta _1+\theta _2,a_2=\theta _1,$$
$$a_0^{}=\theta _1^{},a_1^{}=\theta _0^{}\theta _1^{}+\theta _2^{},a_2^{}=\theta _1^{}.$$
Observe $`a_0=a_2`$ and $`a_0^{}=a_2^{}`$, so that the Leonard pair is balanced. On the other hand, it is essentially bipartite if and only if $`\theta _1=\theta _0\theta _1+\theta _2`$, and it is essentially dual bipartite if and only if $`\theta _1^{}=\theta _0^{}\theta _1^{}+\theta _2^{}`$. Therefore it is not essentially bipartite, and is not essentially dual biparitite for $`2\theta _1\theta _0+\theta _2`$ and $`2\theta _1^{}\theta _0^{}+\theta _2^{}`$.
## 6 Case I: $`d3`$, $`q1`$, $`q1`$
In this section we assume $`d3`$, $`q1`$, $`q1`$.
###### Theorem 6.1
There exist scalars $`\eta `$, $`\mu `$, $`h`$, $`\eta ^{}`$, $`\mu ^{}`$, $`h^{}`$, $`\tau `$ in $`\overline{𝕂}`$ such that for $`0id`$
$`\theta _i`$ $`=`$ $`\eta +\mu q^i+hq^{di},`$ (20)
$`\theta _i^{}`$ $`=`$ $`\eta ^{}+\mu ^{}q^i+h^{}q^{di},`$ (21)
and for $`1id`$
$`\phi _i`$ $`=`$ $`(q^i1)(q^{di+1}1)(\tau \mu \mu ^{}q^{i1}hh^{}q^{di}),`$ (22)
$`\varphi _i`$ $`=`$ $`(q^i1)(q^{di+1}1)(\tau h\mu ^{}q^{i1}\mu h^{}q^{di}).`$ (23)
Proof. These are (27), (28), (31), (32) in after a change of variables. $`\mathrm{}`$
###### Remark 6.2
For $`1id`$ we have $`q^i1`$; otherwise $`\phi _i=0`$ by (22).
###### Lemma 6.3
$`H=(q1)^2((q^{d1}+1)\tau q^{d1}(h+\mu )(h^{}+\mu ^{}))`$.
Proof. It is routine to verify this equation using (4), (20), (21) and (22). $`\mathrm{}`$
###### Lemma 6.4
Assume $`H=0`$. Then $`q^{d1}+10`$ and
$$\tau =\frac{q^{d1}(h+\mu )(h^{}+\mu ^{})}{q^{d1}+1}.$$
(24)
Proof. Assume $`q^{d1}+1=0`$. Then $`1=q^{d1}`$, so that
$`0`$ $`=`$ $`(q^{d1}+1)\tau `$
$`=`$ $`q^{d1}(h+\mu )(h^{}+\mu ^{})`$
$`=`$ $`q^{d1}(\mu hq^{d1})(\mu ^{}h^{}q^{d1})`$
$`=`$ $`q^{d1}(q1)^2(\theta _0\theta _1)(\theta _0^{}\theta _1^{}),`$
a contradiction, so we must have $`q^{d1}+10`$ and (24) follows. $`\mathrm{}`$
###### Lemma 6.5
Assume $`H=0`$. Then the following coincide.
$$\frac{(a_1a_{d1})(\theta _0^{}\theta _3^{})(\theta _{d3}^{}\theta _d^{})}{\theta _0^{}\theta _d^{}},$$
$$\frac{(a_1^{}a_{d1}^{})(\theta _0\theta _3)(\theta _{d3}\theta _d)}{\theta _0\theta _d},$$
$$\frac{(1q^2)(q^31)^2(q^{d1}1)(q^{d2}1)\tau }{q^2(q^d1)}.$$
Proof. It is routine to verify the coincidence using (4), (20), (21), (22) and (24). $`\mathrm{}`$
###### Lemma 6.6
Assume $`H=0`$. Then the following are equivalent.
* $`a_1=a_{d1}`$,
* $`a_1^{}=a_{d1}^{}`$,
* $`(h+\mu )(h^{}+\mu ^{})=0`$,
* $`\tau =0`$.
Proof. Follows from Lemma 6.5 and (24). $`\mathrm{}`$
###### Theorem 6.7
Assume $`d3`$, $`q1`$, $`q1`$. Then the following are equivalent.
* $`a_0=a_d`$ and $`a_1=a_{d1}`$,
* $`a_0^{}=a_d^{}`$ and $`a_1^{}=a_{d1}^{}`$,
* $`a_i=a_{di}`$ and $`a_i^{}=a_{di}^{}`$ for $`0id`$,
* $`\tau =0`$ and $`(h+\mu )(h^{}+\mu ^{})=0`$.
Proof. The conditions (i), (ii), (iv) are equivalent by Lemmas 6.3 and 6.6. Clearly (iii) implies (i). We show (iv) implies (iii). Observe that we have $`h=\mu `$ or $`h^{}=\mu ^{}`$. For the case $`h=\mu `$, it is routine to verify $`a_{di}a_i=0`$ and $`a_{di}^{}a_i^{}=0`$ using (4), (20), (21) and (22) with $`\tau =0`$ and $`h=\mu `$. The case $`h^{}=\mu ^{}`$ is similar. $`\mathrm{}`$
###### Lemma 6.8
For $`0id`$,
$`\theta _i+\theta _{di}`$ $`=`$ $`2\eta +(h+\mu )(q^i+q^{di}),`$
$`\theta _i^{}+\theta _{di}^{}`$ $`=`$ $`2\eta ^{}+(h^{}+\mu ^{})(q^i+q^{di}).`$
Proof. It is routine to verify these equations using (20) and (21). $`\mathrm{}`$
###### Lemma 6.9
For $`1id`$,
$`\phi _i+\varphi _i`$ $`=`$ $`(q^i1)(q^{di+1}1)(2\tau (h+\mu )(\mu ^{}q^{i1}+h^{}q^{di})),`$
$`\phi _i+\varphi _{di+1}`$ $`=`$ $`(q^i1)(q^{di+1}1)(2\tau (h^{}+\mu ^{})(\mu q^{i1}+hq^{di})).`$
Proof. It is routine to verify these equations using (22) and (23). $`\mathrm{}`$
###### Lemma 6.10
The following hold.
* Assume $`\tau =0`$ and $`h^{}+\mu ^{}=0`$. Then
$$a_1a_0=\frac{q^{d2}(q1)(q^21)^2(q^{d1}1)(\mu ^{})^2(h+\mu )}{(\theta _0^{}\theta _1^{})(\theta _1^{}\theta _2^{})}.$$
* Assume $`\tau =0`$ and $`h+\mu =0`$ then
$$a_1^{}a_0^{}=\frac{q^{d2}(q1)(q^21)^2(q^{d1}1)\mu ^2(h^{}+\mu ^{})}{(\theta _0\theta _1)(\theta _1\theta _2)}.$$
Proof. It is routine to verify these equations using (4), (20), (21), (22), (23). $`\mathrm{}`$
###### Theorem 6.11
Assume $`d3`$, $`q1`$, $`q1`$. Then the following are equivalent.
* $`\tau =0`$ and $`h+\mu =0`$.
* $`a_i`$ is independent of $`i`$ for $`0id`$.
* $`\theta _i+\theta _{di}`$ is independent of $`i`$ for $`0id`$, and $`\phi _i=\varphi _i`$ for $`1id`$.
Suppose (i)–(iii) hold. Then the common value of $`a_i`$ is $`\eta `$, and the common value of $`\theta _i+\theta _{di}`$ is $`2\eta `$.
Proof. (i)$``$(ii): Evaluating (20), (22) using $`\tau =0`$ and $`h=\mu `$ we find
$`\theta _i`$ $`=`$ $`\eta +\mu (q^iq^{di})(0id),`$ (25)
$`\phi _i`$ $`=`$ $`\mu (q^i1)(1q^{di+1})(\mu ^{}q^{i1}h^{}q^{di})(1id).`$ (26)
Evaluating the equation on the left in (4) using (21), (25), (26) we routinely find $`a_i=\eta `$ for $`0id`$.
(i)$``$(iii): Setting $`h+\mu =0`$ in Lemma 6.8 we find $`\theta _i+\theta _{di}=2\eta `$ for $`0id`$. Setting $`\tau =0`$ and $`h+\mu =0`$ in Lemma 6.9 we find $`\phi _i=\varphi _i`$ for $`1id`$.
(ii)$``$(i): We have $`\tau =0`$ and $`(h+\mu )(h^{}+\mu ^{})=0`$ by Theorem 6.7. Suppose $`h+\mu 0`$. Then we must have $`h^{}+\mu ^{}=0`$, so that Lemma 6.10 implies $`\mu ^{}(h+\mu )=0`$. Observe that we have $`\mu ^{}0`$; otherwise $`h^{}=h^{}+\mu ^{}=0`$ so that $`\theta _0^{}=\theta _1^{}`$. Hence $`h+\mu =0`$.
(iii)$``$(i): Consider the quantity $`\theta _0+\theta _d\theta _1\theta _{d1}`$. By assumption this quantity is $`0`$. By Lemma 6.8 this quantity is $`(q1)(q^{d1}1)(h+\mu )`$ so $`h+\mu =0`$. Setting $`\phi _i+\varphi _i=0`$, $`h+\mu =0`$ in Lemma 6.9 we find $`2\tau =0`$. Observe $`\text{Char}(𝕂)2`$; otherwise $`\theta _d=\theta _0`$ by Lemma 6.8. We conclude $`\tau =0`$. $`\mathrm{}`$
###### Theorem 6.12
Assume $`d3`$, $`q1`$, $`q1`$. Then the following are equivalent.
* $`\tau =0`$ and $`h^{}+\mu ^{}=0`$.
* $`a_i^{}`$ is independent of $`i`$ for $`0id`$.
* $`\theta _i^{}+\theta _{di}^{}`$ is independent of $`i`$ for $`0id`$, and $`\phi _i=\varphi _{di+1}`$ for $`1id`$.
Suppose (i)–(iii) hold. Then the common value of $`a_i^{}`$ is $`\eta ^{}`$, and the common value of $`\theta _i^{}+\theta _{di}^{}`$ is $`2\eta ^{}`$.
Proof. Similar to the proof of Thoerem 6.11. $`\mathrm{}`$
## 7 Case II: $`d3`$, $`q=1`$, $`\text{Char}(𝕂)2`$
In this section we assume $`d3`$, $`q=1`$, $`\text{Char}(𝕂)2`$.
###### Theorem 7.1
There exist scalars $`\eta `$, $`\mu `$, $`h`$, $`\eta ^{}`$, $`\mu ^{}`$, $`h^{}`$, $`\tau `$ in $`\overline{𝕂}`$ such that for $`0id`$
$`\theta _i`$ $`=`$ $`\eta +\mu (id/2)+hi(di),`$ (27)
$`\theta _i^{}`$ $`=`$ $`\eta ^{}+\mu ^{}(id/2)+h^{}i(di),`$ (28)
and for $`1id`$
$`\phi _i`$ $`=`$ $`i(di+1)(\tau \mu \mu ^{}/2+(h\mu ^{}+\mu h^{})(i(d+1)/2)+hh^{}(i1)(di)),`$ (29)
$`\varphi _i`$ $`=`$ $`i(di+1)(\tau +\mu \mu ^{}/2+(h\mu ^{}\mu h^{})(i(d+1)/2)+hh^{}(i1)(di)).`$ (30)
Proof. These are (35), (36), (38), (39) in after a change of variables. $`\mathrm{}`$
###### Remark 7.2
If $`h=0`$ then $`\mu 0`$, otherewise $`\theta _1=\theta _0`$. Similarly if $`h^{}=0`$ then $`\mu ^{}0`$. For any prime $`i`$ such that $`id`$ we have $`\text{Char}(𝕂)i`$; otherwise $`\phi _i=0`$ by (29).
###### Lemma 7.3
$`H=2\tau +hh^{}(d1)^2`$.
Proof. It is routine to verify this equation using (4), (27), (28) and (29). $`\mathrm{}`$
###### Lemma 7.4
Assume $`H=0`$. Then
$$\tau =hh^{}(d1)^2/2.$$
(31)
Proof. Follows from Lemma 7.3. $`\mathrm{}`$
###### Lemma 7.5
Assume $`H=0`$. Then the following coincide.
$$\frac{(a_1a_{d1})(\theta _0^{}\theta _3^{})(\theta _{d3}^{}\theta _d^{})}{\theta _0^{}\theta _d^{}},$$
$$\frac{(a_1^{}a_{d1}^{})(\theta _0\theta _3)(\theta _{d3}\theta _d)}{\theta _0\theta _d},$$
$$36d^1(d1)(d2)hh^{}.$$
Proof. It is routine to verify the coincidence using (4), (27), (28), (29) and (31). $`\mathrm{}`$
###### Lemma 7.6
Assume $`H=0`$. Then the following are equivalent.
* $`a_1=a_{d1}`$,
* $`a_1^{}=a_{d1}^{}`$,
* $`hh^{}=0`$,
* $`\tau =0`$.
Proof. Follows from Lemma 7.5 and Remark 7.2. $`\mathrm{}`$
###### Theorem 7.7
Assume $`d3`$, $`q=1`$, $`\text{Char}(𝕂)2`$. Then the following are equivalent.
* $`a_0=a_d`$ and $`a_1=a_{d1}`$,
* $`a_0^{}=a_d^{}`$ and $`a_1^{}=a_{d1}^{}`$,
* $`a_i=a_{di}`$ and $`a_i^{}=a_{di}^{}`$ for $`0id`$,
* $`hh^{}=0`$ and $`\tau =0`$.
Proof. The conditions (i), (ii), (iv) are equivalent by Lemmas 7.3 and 7.6. Clearly (iii) implies (i). We show (iv) implies (iii). It is routine to verify $`a_ia_{di}=0`$ and $`a_i^{}a_{di}^{}=0`$ for each case of $`h=0`$, $`h^{}=0`$ by using (4), (27), (28), (29) with $`\tau =0`$. $`\mathrm{}`$
###### Lemma 7.8
For $`0id`$,
$`\theta _i+\theta _{di}`$ $`=`$ $`2(\eta +hi(di)),`$
$`\theta _i^{}+\theta _{di}^{}`$ $`=`$ $`2(\eta ^{}+h^{}i(di)).`$
Proof. It is routine to verify these equations using (27) and (28). $`\mathrm{}`$
###### Lemma 7.9
For $`1id`$,
$`\phi _i+\varphi _i`$ $`=`$ $`i(di+1)(2\tau (d2i+1)h\mu ^{}+2hh^{}(di)(i1)),`$
$`\phi _i+\varphi _{di+1}`$ $`=`$ $`i(di+1)(2\tau (d2i+1)h^{}\mu +2hh^{}(di)(i1)).`$
Proof. It is routine to verify these equations using (29) and (30). $`\mathrm{}`$
###### Lemma 7.10
The following hold.
* Assume $`\tau =0`$ and $`h^{}=0`$. Then
$$a_0a_1=2(d1)h.$$
* Assume $`\tau =0`$ and $`h=0`$. Then
$$a_0^{}a_1^{}=2(d1)h^{}.$$
Proof. It is routine to verify these equations using (4), (27), (28), (29). $`\mathrm{}`$
###### Theorem 7.11
Assume $`d3`$, $`q=1`$, $`\text{Char}(𝕂)2`$. Then the following are equivalent.
* $`h=0`$ and $`\tau =0`$.
* $`a_i`$ is independent of $`i`$ for $`0id`$.
* $`\theta _i+\theta _{di}`$ is independent of $`i`$ for $`0id`$, and $`\phi _i=\varphi _i`$ for $`1id`$.
Suppose (i)–(iii) hold. Then the common value of $`a_i`$ is $`\eta `$, and the common value of $`\theta _i+\theta _{di}`$ is $`2\eta `$.
Proof. (i)$``$(ii): Evaluating (27), (29) using $`h=0`$ and $`\tau =0`$ we find
$`\theta _i`$ $`=`$ $`\eta +(id/2)\mu (0id),`$ (32)
$`\phi _i`$ $`=`$ $`i(di+1)\mu (\mu ^{}+h^{}(d2i+1))/2(1id).`$ (33)
Evaluating the equation on the left in (4) using (28), (32), (33) we routinely find $`a_i=\eta `$ for $`0id`$.
(i)$``$(iii): Setting $`h=0`$ in Lemma 7.8 we find $`\theta _i+\theta _{di}=2\eta `$ for $`0id`$. Setting $`h=0`$ and $`\tau =0`$ in Lemma 7.9 we find $`\phi _i=\varphi _i`$ for $`1id`$.
(ii)$``$(i): We have $`\tau =0`$ and $`hh^{}=0`$ by Theorem 7.7. Suppose $`h0`$. Then we must have $`h^{}=0`$. Then Lemma 7.10 implies $`h=0`$.
(iii)$``$(i): Consider the quantity $`\theta _0+\theta _d\theta _1\theta _{d1}`$. By assumption this quantity is $`0`$. By Lemma 7.8 this quantity is $`2(1d)h`$ so $`h=0`$. Setting $`\phi _i+\varphi _i=0`$, $`h=0`$ in Lemma 7.9 we find $`\tau =0`$. $`\mathrm{}`$
###### Theorem 7.12
Assume $`d3`$, $`q=1`$, $`\text{Char}(𝕂)2`$. Then the following are equivalent.
* $`h^{}=0`$ and $`\tau =0`$.
* $`a_i^{}`$ is independent of $`i`$ for $`0id`$.
* $`\theta _i^{}+\theta _{di}^{}`$ is independent of $`i`$ for $`0id`$, and $`\phi _i=\varphi _{di+1}`$ for $`1id`$.
Suppose (i)–(iii) hold. Then the common value of $`a_i^{}`$ is $`\eta ^{}`$, and the common value of $`\theta _i^{}+\theta _{di}^{}`$ is $`2\eta ^{}`$.
Proof. Similar to the proof of Theorem 7.11. $`\mathrm{}`$
## 8 Case III: $`d3`$, $`q=1`$, $`\text{Char}(𝕂)2`$, $`d`$ even
In this section we assume $`d3`$, $`q=1`$, $`\text{Char}(𝕂)2`$, and $`d`$ is even.
###### Theorem 8.1
\[4, Theorem 5.16, Example 5.14\] There exist scalars $`\eta `$, $`h`$, $`s`$, $`\eta ^{}`$, $`h^{}`$, $`s^{}`$, $`\tau `$ in $`\overline{𝕂}`$ such that for $`0id`$
$`\theta _i`$ $`=`$ $`\{\begin{array}{cc}\eta +s+h(id/2)\hfill & \text{if }i\text{ is even},\hfill \\ \eta sh(id/2)\hfill & \text{if }i\text{ is odd},\hfill \end{array}`$ (34)
$`\theta _i^{}`$ $`=`$ $`\{\begin{array}{cc}\eta ^{}+s^{}+h^{}(id/2)\hfill & \text{if }i\text{ is even},\hfill \\ \eta ^{}s^{}h^{}(id/2)\hfill & \text{if }i\text{ is odd},\hfill \end{array}`$ (35)
and for $`1id`$
$`\phi _i`$ $`=`$ $`\{\begin{array}{cc}i(\tau sh^{}s^{}hhh^{}(i(d+1)/2))\hfill & \text{if }i\text{ is even},\hfill \\ (di+1)(\tau +sh^{}+s^{}h+hh^{}(i(d+1)/2))\hfill & \text{if }i\text{ is odd},\hfill \end{array}`$ (36)
$`\varphi _i`$ $`=`$ $`\{\begin{array}{cc}i(\tau sh^{}+s^{}h+hh^{}(i(d+1)/2))\hfill & \text{if }i\text{ is even},\hfill \\ (di+1)(\tau +sh^{}s^{}hhh^{}(i(d+1)/2))\hfill & \text{if }i\text{ is odd}.\hfill \end{array}`$ (37)
Proof. These are (19)–(22) in after a change of variables. $`\mathrm{}`$
###### Remark 8.2
We have $`h0`$; otherwise $`\theta _0=\theta _2`$ by (34). Similary we have $`h^{}0`$. For any prime $`i`$ such that $`id/2`$ we have $`\text{Char}(𝕂)i`$; otherwise $`\phi _{2i}=0`$ by (36). By this and since $`\text{Char}(𝕂)2`$ we find $`\text{Char}(𝕂)`$ is either $`0`$ or an odd prime greater than $`d/2`$. Observe neither of $`d`$, $`d2`$ vanish in $`𝕂`$ since otherwise $`\text{Char}(𝕂)`$ must divide $`d/2`$ or $`(d2)/2`$.
###### Lemma 8.3
$`H=2(d1)\tau +4ss^{}`$.
Proof. It is routine to verify this equation using (4), (34), (35), (36). $`\mathrm{}`$
###### Lemma 8.4
Assume $`H=0`$. Then $`d1`$ is nonzero in $`𝕂`$ and
$$\tau =\frac{2ss^{}}{1d}.$$
(38)
Proof. Suppose $`d1`$ is zero in $`𝕂`$. Then Lemma 8.3 implies $`ss^{}=0`$. If $`s=0`$ then $`\theta _1=\theta _0`$ by (34). If $`s^{}=0`$ then $`\theta _1^{}=\theta _0^{}`$ by (35). Hence $`d1`$ is nonzero and (38) follows. $`\mathrm{}`$
###### Lemma 8.5
Assume $`H=0`$. Then the following coincide.
$$\frac{(a_1a_{d1})(\theta _0^{}\theta _3^{})(\theta _{d3}^{}\theta _d^{})}{\theta _0^{}\theta _d^{}},$$
$$\frac{(a_1^{}a_{d1}^{})(\theta _0\theta _3)(\theta _{d3}\theta _d)}{\theta _0\theta _d},$$
$$\frac{16(d2)ss^{}}{d(d1)}.$$
Proof. It is routine to verify the coincidence using (4), (34), (35), (36) and (38). $`\mathrm{}`$
###### Lemma 8.6
Assume $`H=0`$. Then the following are equivalent.
* $`a_1=a_{d1}`$,
* $`a_1^{}=a_{d1}^{}`$,
* $`ss^{}=0`$,
* $`\tau =0`$.
Proof. Follows from Lemma 8.5 and (38). $`\mathrm{}`$
###### Theorem 8.7
Assume $`d3`$, $`q=1`$, $`\text{Char}(𝕂)2`$, and $`d`$ is even. Then the following are equivalent.
* $`a_0=a_d`$ and $`a_1=a_{d1}`$,
* $`a_0^{}=a_d^{}`$ and $`a_1^{}=a_{d1}^{}`$,
* $`a_i=a_{di}`$ and $`a_i^{}=a_{di}^{}`$ for $`0id`$,
* $`ss^{}=0`$ and $`\tau =0`$.
Proof. The conditions (i), (ii), (iv) are equivalent by Lemmas 8.3 and 8.6. Clearly (iii) implies (i). We show (iv) implies (iii). It is routine to verify $`a_ia_{di}=0`$ and $`a_i^{}a_{di}^{}=0`$ for each case of $`s=0`$ and $`s^{}=0`$ using (4), (34), (35), (36) with $`\tau =0`$. $`\mathrm{}`$
###### Lemma 8.8
For $`0id`$,
$`\theta _i+\theta _{di}`$ $`=`$ $`\{\begin{array}{cc}2(\eta +s)\hfill & \text{if }i\text{ is even},\hfill \\ 2(\eta s)\hfill & \text{if }i\text{ is odd},\hfill \end{array}`$
$`\theta _i^{}+\theta _{di}^{}`$ $`=`$ $`\{\begin{array}{cc}2(\eta ^{}+s^{})\hfill & \text{if }i\text{ is even},\hfill \\ 2(\eta ^{}s^{})\hfill & \text{if }i\text{ is odd}.\hfill \end{array}`$
Proof. It is routine to verify these equations using (34) and (35). $`\mathrm{}`$
###### Lemma 8.9
For $`1id`$,
$`\phi _i+\varphi _i`$ $`=`$ $`\{\begin{array}{cc}2i(\tau sh^{})\hfill & \text{if }i\text{ is even},\hfill \\ 2(di+1)(\tau +sh^{})\hfill & \text{if }i\text{ is odd},\hfill \end{array}`$
$`\phi _i+\varphi _{di+1}`$ $`=`$ $`\{\begin{array}{cc}2i(\tau s^{}h)\hfill & \text{if }i\text{ is even},\hfill \\ 2(di+1)(\tau +s^{}h)\hfill & \text{if }i\text{ is odd}.\hfill \end{array}`$
Proof. It is routine to verify these equations using (36) and (37). $`\mathrm{}`$
###### Lemma 8.10
The following hold.
* Assume $`\tau =0`$ and $`s^{}=0`$. Then each of $`d1`$, $`d3`$ is nonzero in $`𝕂`$ and
$$a_0a_1=\frac{4s}{(d1)(d3)}.$$
* Assume $`\tau =0`$ and $`s=0`$. Then each of $`d1`$, $`d3`$ is nonzero in $`𝕂`$ and
$$a_0^{}a_1^{}=\frac{4s^{}}{(d1)(d3)}.$$
Proof. We show (i). If $`d1`$ is zero in $`𝕂`$ then $`\theta _0^{}=\theta _1^{}`$ by (35). If $`d3`$ is zero in $`𝕂`$ then $`\theta _0^{}=\theta _3^{}`$ by (35). Hence each of $`d1`$, $`d3`$ is nonzero. Now we routinely find the equation for $`a_0a_1`$ using (4), (34), (35) and (36). The proof of (ii) is similar. $`\mathrm{}`$
###### Theorem 8.11
Assume $`d3`$, $`q=1`$, $`\text{Char}(𝕂)2`$, and $`d`$ is even. Then the following are equivalent.
* $`s=0`$ and $`\tau =0`$.
* $`a_i`$ is independent of $`i`$ for $`0id`$.
* $`\theta _i+\theta _{di}`$ is independent of $`i`$ for $`0id`$, and $`\phi _i=\varphi _i`$ for $`1id`$.
Suppose (i)–(iii) hold. Then the common value of $`a_i`$ is $`\eta `$, and the common value of $`\theta _i+\theta _{di}`$ is $`2\eta `$.
Proof. (i)$``$(ii): Evaluating (34), (36) using $`s=0`$ and $`\tau =0`$ we find
$`\theta _i`$ $`=`$ $`\{\begin{array}{cc}\eta +h(id/2)\hfill & \text{ if }i\text{ is even},\hfill \\ \eta h(id/2)\hfill & \text{ if }i\text{ is odd},\hfill \end{array}`$ (39)
$`\phi _i`$ $`=`$ $`\{\begin{array}{cc}hi(s^{}+h^{}(i(d+1)/2))\hfill & \text{ if }i\text{ is even},\hfill \\ h(di+1)(s^{}+h^{}(i(d+1)/2))\hfill & \text{ if }i\text{ is odd}.\hfill \end{array}`$ (40)
Evaluating the equation on the left in (4) using (39), (40) we routinely find $`a_i=\eta `$ for $`0id`$.
(i)$``$(iii): Setting $`s=0`$ in Lemma 8.8 we find $`\theta _i+\theta _{di}=2\eta `$ for $`0id`$. Setting $`s=0`$ and $`\tau =0`$ in Lemma 8.9 we find $`\phi _i=\varphi _i`$ for $`1id`$.
(ii)$``$(i): Suppose (i) does not hold. Then from Theorem 8.7, we must have $`s^{}=0`$ and $`\tau =0`$. From our assuption, we have $`a_1a_0=0`$, so Lemma 8.10 implies $`s=0`$, a contradiction.
(iii)$``$(i): From Lemma 8.9 for $`i=1,2`$,
$$0=\phi _1+\varphi _1=2d(\tau +sh^{}),$$
$$0=\phi _2+\varphi _2=4(\tau sh^{}).$$
These equations imply $`\tau =0`$ and $`s=0`$. $`\mathrm{}`$
###### Theorem 8.12
Assume $`d3`$, $`q=1`$, $`\text{Char}(𝕂)2`$, and $`d`$ is even. Then the following are equivalent.
* $`s^{}=0`$ and $`\tau =0`$.
* $`a_i^{}`$ is independent of $`i`$ for $`0id`$.
* $`\theta _i^{}+\theta _{di}^{}`$ is independent of $`i`$ for $`0id`$, and $`\phi _i=\varphi _{di+1}`$ for $`1id`$.
Suppose (i)–(iii) hold. Then the common value of $`a_i^{}`$ is $`\eta ^{}`$, and the common value of $`\theta _i^{}+\theta _{di}^{}`$ is $`2\eta ^{}`$.
Proof. Similar to the proof of Theorem 8.11. $`\mathrm{}`$
## 9 Case IV: $`d3`$, $`q=1`$, $`\text{Char}(𝕂)2`$, $`d`$ odd
In this section we assume $`d3`$, $`q=1`$, $`\text{Char}(𝕂)2`$, and $`d`$ is odd.
###### Theorem 9.1
\[4, Theorem 5.16, Example 5.14\] There exist scalars $`\eta `$, $`h`$, $`s`$, $`\eta ^{}`$, $`h^{}`$, $`s^{}`$, $`\tau `$ in $`\overline{𝕂}`$ such that for $`0id`$
$`\theta _i`$ $`=`$ $`\{\begin{array}{cc}\eta +s+h(id/2)\hfill & \text{if }i\text{ is even},\hfill \\ \eta sh(id/2)\hfill & \text{if }i\text{ is odd},\hfill \end{array}`$ (41)
$`\theta _i^{}`$ $`=`$ $`\{\begin{array}{cc}\eta ^{}+s^{}+h^{}(id/2)\hfill & \text{if }i\text{ is even},\hfill \\ \eta ^{}s^{}h^{}(id/2)\hfill & \text{if }i\text{ is odd},\hfill \end{array}`$ (42)
and for $`1id`$
$`\phi _i`$ $`=`$ $`\{\begin{array}{cc}hh^{}i(di+1)\hfill & \text{if }i\text{ is even},\hfill \\ \tau 2ss^{}+i(di+1)hh^{}2(hs^{}+h^{}s)(i(d+1)/2)\hfill & \text{if }i\text{ is odd},\hfill \end{array}`$ (43)
$`\varphi _i`$ $`=`$ $`\{\begin{array}{cc}hh^{}i(di+1)\hfill & \text{if }i\text{ is even},\hfill \\ \tau +2ss^{}+i(di+1)hh^{}2(hs^{}h^{}s)(i(d+1)/2)\hfill & \text{if }i\text{ is odd}.\hfill \end{array}`$ (44)
###### Remark 9.2
Observe $`hh^{}0`$, and $`\text{Char}(𝕂)`$ is either $`0`$ or an odd prime greater than $`d/2`$. Also observe $`d1`$ does not vanish in $`𝕂`$. These can be observed in a similar way as Remark 8.2.
###### Lemma 9.3
$`H=2\tau +(d^2+1)hh^{}`$.
Proof. It is routine to verify this equation using (4), (41), (42) and (43). $`\mathrm{}`$
###### Lemma 9.4
Assume $`H=0`$. Then
$$\tau =(d^2+1)hh^{}/2.$$
(45)
Proof. Follows from Lemma 9.3. $`\mathrm{}`$
###### Lemma 9.5
Assume $`H=0`$. Then the following coincide.
$$\frac{(a_1a_{d1})(\theta _0^{}\theta _3^{})(\theta _{d3}^{}\theta _d^{})}{\theta _0^{}\theta _d^{}},$$
$$\frac{(a_1^{}a_{d1}^{})(\theta _0\theta _3)(\theta _{d3}\theta _d)}{\theta _0\theta _d},$$
$$4(d1)hh^{}.$$
Proof. It is routine to verify the coincidence using (4), (41), (42), (43) and (45). $`\mathrm{}`$
###### Theorem 9.6
Assume $`d3`$, $`q=1`$, $`\text{Char}(𝕂)2`$, and $`d`$ is odd. If $`a_0=a_d`$ then $`a_1a_{d1}`$ and $`a_1^{}a_{d1}^{}`$.
Proof. Follows from Lemma 9.5 and Remark 9.2. $`\mathrm{}`$
###### Theorem 9.7
Assume $`d3`$, $`q=1`$, $`\text{Char}(𝕂)2`$, and $`d`$ is odd. Then $`\phi _2+\varphi _20`$ and $`\phi _2+\varphi _{d1}0`$.
Proof. From (43) and (44), $`\phi _2+\varphi _2=\phi _2+\varphi _{d1}=4(d1)hh^{}0`$ by Remark 9.2. $`\mathrm{}`$
## 10 Case V: $`d3`$, $`q=1`$, $`\text{Char}(𝕂)=2`$
In this section we assume $`d3`$, $`q=1`$, $`\text{Char}(𝕂)=2`$.
###### Theorem 10.1
\[4, Theorem 5.16, Example 5.15\] We have $`d=3`$, and there exist scalars $`h`$, $`s`$, $`h^{}`$, $`s^{}`$, $`r`$ in $`\overline{𝕂}`$ such that
$$\begin{array}{ccc}\theta _1=\theta _0+h(s+1),\hfill & \theta _2=\theta _0+h,\hfill & \theta _3=\theta _0+hs,\hfill \\ \theta _1^{}=\theta _0^{}+h^{}(s^{}+1),\hfill & \theta _2^{}=\theta _0^{}+h^{},\hfill & \theta _3^{}=\theta _0^{}+h^{}s^{},\hfill \\ \phi _1=hh^{}r,\hfill & \phi _2=hh^{},\hfill & \phi _3=hh^{}(r+s+s^{}),\hfill \\ \varphi _1=hh^{}(r+s(1+s^{})),\hfill & \varphi _2=hh^{},\hfill & \varphi _3=hh^{}(r+s^{}(1+s)).\hfill \end{array}$$
###### Remark 10.2
Each of $`h`$, $`h^{}`$, $`s`$, $`s^{}`$ is nonzero, and each of $`s`$, $`s^{}`$ is not equal to $`1`$.
###### Lemma 10.3
$`a_0a_3`$ $`=`$ $`{\displaystyle \frac{hs^{}(1+s)}{1+s^{}}},`$ (46)
$`a_0^{}a_3^{}`$ $`=`$ $`{\displaystyle \frac{h^{}s(1+s^{})}{1+s}}.`$ (47)
Proof. Obtained by a routine computation. We remark that $`2=0`$ and $`1=1`$ since $`\text{Char}(𝕂)=2`$. $`\mathrm{}`$
###### Theorem 10.4
Assume $`d3`$, $`q=1`$, $`\text{Char}(𝕂)=2`$. Then $`a_0a_d`$ and $`a_0^{}a_d^{}`$.
Proof. Immediate from Lemma 10.3 and since none of $`h`$, $`h^{}`$, $`s`$, $`s^{}`$, $`1+s`$, $`1+s^{}`$ is zero. $`\mathrm{}`$
###### Theorem 10.5
Assume $`d3`$, $`q=1`$, $`\text{Char}(𝕂)=2`$. Then $`\phi _1+\varphi _10`$ and $`\phi _1+\varphi _d0`$.
Proof. We have
$`\phi _1+\varphi _1`$ $`=`$ $`hh^{}s(1+s^{}),`$
$`\phi _1+\varphi _3`$ $`=`$ $`hh^{}s^{}(1+s).`$
These values are nonzero since none of $`h`$, $`h^{}`$, $`s`$, $`s^{}`$, $`1+s`$, $`1+s^{}`$ is zero. $`\mathrm{}`$
Kazumasa Nomura
College of Liberal Arts and Sciences
Tokyo Medical and Dental University
Kohnodai, Ichikawa, 272-0827 Japan
email: nomura.las@tmd.ac.jp
Paul Terwilliger
Department of Mathematics
University of Wisconsin
480 Lincoln drive, Madison, Wisconsin, 53706 USA
email: terwilli@math.wisc.edu
Keywords. Leonard pair, Terwilliger algebra, Askey scheme, $`q`$-Racah polynomial.
2000 Mathematics Subject Classification. 05E30, 05E35, 33C45, 33D45. |
warning/0506/math0506164.html | ar5iv | text | # Integrable Systems and Harmonic Maps into Lie Groups
## I Introduction
The concept of integrability in infinite dimensions is not clear cut. In classical mechanics, a Hamiltonian system with a $`2n`$-dimensional phase space is said to be integrable if it has $`n`$ constants of the motion in involution (with vanishing Poisson brackets). In systems of partial diferential equations, there are infinitely many degrees of freedom, and there is no straight-forward corresponding definition. Since the equations we are working with are static equations, not flow equations, integrals are not even defined.
There are other features of the systems that are associated with being ‘integrable’. First, the equations are to some degree soluble, meaning that explicit solutions can be found and there exist general methods for constructing solutions, which may be superimposable in some extended sense. Alternatively, it may be possible to find a large number of constants of the motion, or the system may have the Painleve property.
The prototype of the integrable system is the Korteweg-De Vries (KdV) equation. A solitary wave is modelled by a soliton solution to the KdV equation:
$$4u_tu_{xxx}6uu_x=0$$
(1)
where $`\kappa `$ and $`x_0`$ are constants. The integrability of the KdV equation can be traced to the existence of a Lax pair, or obtaining the equation is the condition that the two differential operators commute. For KdV, the equation is the condition that the two differential operators
$$L=_x^2+uandM=_t_x^3\frac{3}{2}u_x\frac{3}{4}u_x$$
(2)
commute.
In this paper we look at the nature of relationships between some of the better-known integrable systems. We focus on one of the most widely studied integrable system, the Toda model, and find a direct connection to the theory of harmonic maps into Lie groups.
## II Toda Theory
The Toda field is a multicomponent field in two dimensions satisfying
$$^2\varphi _i+k_{ij}e^{\varphi _j}=0$$
(3)
where $`k_{ij}`$ is the Cartan matrix of a semisimple Lie algebra $``$. The Liouville equation is obtained from (3) by taking $`k`$ as the Cartan matrix of the Lie algebra sl(2,C). The Toda field equations have a formulation in terms of a Lax pair. First, we need some results from the theory of semisimple Lie algebra to describe this. If a Lie algebra $``$ has a basis $`L^i`$ and structure constants $`f^ij_k`$ then
$$[L^i,L^j]=\underset{k}{}f_k^{ij}L_k$$
(4)
Associated with $``$ is its set of roots $`\mathrm{\Phi }`$. These roots are $`r`$-dimensional vectors, where the rank $`r`$ of $``$ is the maximal number of linearly independent commuting generators. A set of simple roots $`\mathrm{\Delta }=\alpha _i;i=1\mathrm{}r`$ is a subset of $`\mathrm{\Phi }`$ such that the difference of any two of its elements is not in $`\mathrm{\Phi }`$. Any root may be expressed
$$\alpha =\underset{i}{}n_i\alpha _i$$
(5)
with the $`n_i`$ either all non-negative or all non-positive integers. Hence $`\mathrm{\Phi }`$ is divided into two sets $`\mathrm{\Phi }^+`$ and $`\mathrm{\Phi }^{}`$ containing positive and negative roots respectively.
The Cartan subalgebra is the maximal abelian subalgebra of $``$ is called the Cartan subalgebra. For $`SU(n)`$, these are just the diagonal matrices. A Chevalley basis of the Lie algebra is a choice of an element $`H_i`$ of the Cartan for each simple root. The remaining generators $`E_\alpha `$ are labelled by roots. Those labelled by simple roots are for convenience denoted as $`E_i^\pm :=E_{\pm \alpha _i}`$, which for our purposes are just the generators for the off-diagonal parts of the algebra.
In this basis the algebra takes the form
$$\begin{array}{c}[H_i,H_j]=0\\ [H_i,E_j^\pm ]=\pm k_{ji}E_j^\pm \\ [E_i^+,E_j^{}]=\delta _{ij}H_j\end{array}$$
(6)
The upper and lower triangular matrices are given by
$$_\pm =\{H_i,E_\alpha :i=1\mathrm{}r,\alpha \mathrm{\Phi }^\pm \}$$
(7)
We are now ready to write down the Lax pair and the associated linear problem which codify the Toda theory.
We follow Leznov and Saveliev and consider a connection whose components have values in different subalgebras
$$A_{z,\overline{z}}_{+,}$$
(8)
If the connection is trivialisable, that is, it satisfies the linear equations
$$_zg=A_zg,_{\overline{z}}g=A_{\overline{z}}g$$
(9)
then it automatically satisfies the flatness or integrability condition
$$[_z+A_z,_{\overline{z}}+A_{\overline{z}}]=0$$
(10)
An appropriate choice for $`A`$ is:
$$\begin{array}{c}A_z=\underset{i}{}\left(_z\psi _{}H_i+\alpha E_i^+\right)\\ A_{\overline{z}}=\underset{i}{}\alpha e^{\beta \varphi _i}E_i^{}\end{array}$$
(11)
where $`\alpha ,\beta `$ are constants and $`\psi ,\varphi `$ are two $`r`$-component classical fields. The algebra is straightforward, and in characteristic fashion we get an equation for each element of the Chevalley basis that appears. Write out (10) explicitly, to get
$$\begin{array}{c}_zA_{\overline{z}}_{\overline{z}}A_z+\underset{ij}{}\left(_+\psi _i\alpha e^{\beta \varphi _j}[H_i,E_j^{}]+\alpha ^2e^{\beta \varphi _j}[E_i^+,E_j^{}]\right)\\ =\underset{i}{}\left(\alpha \beta _+\varphi _ie^{\beta \varphi _i}E_i^{}_z_{\overline{z}}\psi _iH_i\right)\underset{ij}{}\alpha _+\psi _ie^{\beta \varphi _j}k_{ji}E_j^{}+\underset{i}{}\alpha ^2e^{\beta \varphi _i}H_i\\ =0\end{array}$$
(12)
which is satisfied by
$$\begin{array}{c}\psi _i=\beta \underset{j}{}k_{ij}^1\varphi _j\\ _z_{\overline{z}}\psi _i=\alpha ^2e^{\beta \varphi _i}\end{array}$$
(13)
or written in the form of a single equation, the Toda equation itself
$$_z_{\overline{z}}\varphi _i=\frac{\alpha ^2}{\beta }\underset{j}{}k_{ij}e^{\beta \varphi _j}$$
(14)
## III Conformal Affine Toda Theory
Because of the appearance of the spectral parameter $`\lambda `$ in our Lax pair for the harmonic map, we will find it is useful to study Lie algebras containing an affine parameter, i.e., affine Lie algebras. The corresponding Toda theories are the Affine Toda theories. We will be able to describe a relationship between harmonic maps into Lie groups and all of these models, most interestingly with the so-called Conformal Affine Toda theory. This conformally invariant field theory is based on the affine Lie algebra $`\widehat{sl_2}`$ which reduces under certain circumstances to the Liouville theory and the non-conformal sinh-Gordon theory.
We first describe the conformally invariant CAT theory . These models are obtained from the usual Toda field theory by adding two fields which transform in the correct way under conformal transformations. The addition of these fields is facilitated by constructing the affine Lie algebra $`\widehat{sl_2}`$. This is the Lie algebra of traceless 2x2 matrices with entries which are Laurent polynomials in $`\lambda `$ (the loop algebra $`\stackrel{~}{sl_2}`$). This algebra is centrally extended as follows:
$$\stackrel{}{sl_2^{}}=\stackrel{}{sl_2}𝒞c.$$
(15)
where
$$[\widehat{X},\widehat{Y}]_{}=[\stackrel{~}{X},\stackrel{~}{Y}]_{}+\frac{1}{2i\pi }𝑑\lambda tr\left[_\lambda \stackrel{~}{X}\left(\lambda \right)\stackrel{~}{Y}\left(\lambda \right)\right]c$$
(16)
The affine Lie algebra $`\widehat{sl_2}`$ is obtained by adding the derivation $`d=\lambda \frac{d}{d\lambda }`$. The algebra can be decomposed as
$$\stackrel{}{sl_2^{}}=𝒩_{}𝒩_+$$
(17)
where $`𝒩_{}`$$`𝒩+`$ are lower and upper triangular matrices respectively, and $``$, the Cartan sub-algebra, is spanned by the elements H, c and d. We write in co-ordinates $`\mathrm{\Phi }:𝒞`$
$$\mathrm{\Phi }=\frac{1}{2}\varphi H+\eta d+\frac{1}{2}\xi c$$
(18)
where $`H,d,c`$ generate $``$.
We follow the construction of Toda field theory shown above, but using complex co-ordinates and a slightly different connection. The Lax pair $`(_z+𝒜_z,_{\overline{z}}+𝒜_{\overline{z}})`$ is written as
$`A_z`$ $`=`$ $`_z\mathrm{\Phi }+e^\mathrm{\Phi }(E_++\lambda E_{})e^\mathrm{\Phi }`$ (19)
$`A_{\overline{z}}`$ $`=`$ $`_{\overline{z}}\mathrm{\Phi }+e^\mathrm{\Phi }(E_{}+\lambda ^1E_+)e^\mathrm{\Phi }.`$ (20)
The zero curvature condition (10) now gives us, after a little effort, the following set of equations
$`_z_{\overline{z}}\varphi `$ $`=`$ $`e^{2\varphi }e^{2\eta 2\varphi }`$ (21)
$`_z_{\overline{z}}\eta `$ $`=`$ $`0`$ (22)
$`_z_{\overline{z}}\xi `$ $`=`$ $`e^{2\eta 2\varphi }.`$ (23)
This system of three equations is the Conformal Affine Toda theory. The reduction to the sinh-Gordon and Liouville theories are realised by the limits $`\eta 0`$ and $`\eta \mathrm{}`$ respectively.
The set of equations (90)-(92) are conformally invariant, as is the reduction to the Liouville equation.
## IV Wess-Zumino-Witten models and Liouville theory
O’Raifeartaigh et al have shown that Liouville theory can be regarded as a reduced SL(2,R) Wess-Zumino-Witten theory. Their reduction uses a decomposition into local fields.
The energy or action for a group-valued field $`g`$ is
$$S(g)=\frac{k}{8\pi }d^2zTr\left[(g^1_zg)(g^1_{\overline{z}}g)\right]+\frac{k}{12\pi }\underset{B^3}{}Tr\left[(g^1dg)^3\right]$$
(24)
Any connected semi-simple real Lie group $`G`$ admits a Gauss decomposition
$$G=XYZ$$
(25)
where $`Y`$ is the direct product $`Y=AK`$ of a simply-connected abelian group $`A`$ and a connected semisimple compact group $`K`$, and the groups $`X`$ and $`Z`$ are simply connected and nilpotent; two arbitrary decompositions are connected by an automorphism of $`G`$. \[BR\]
In the case of $`SL(2,R)`$, there exists such a decomposition for regular g which valid in a neighbourhood of the identity, as follows:
$$g=ABC$$
(26)
where
$`A=\left(\begin{array}{cc}1& x\\ 0& 1\end{array}\right)=\mathrm{exp}(xE_+)`$ (27)
$`C=\left(\begin{array}{cc}1& 0\\ y& 1\end{array}\right)=\mathrm{exp}(yE_{})`$ (28)
$`B=\left(\begin{array}{cc}\mathrm{exp}(\frac{1}{2}\varphi )& 0\\ 0& \mathrm{exp}(\frac{1}{2}\varphi )\end{array}\right)=\mathrm{exp}(\frac{1}{2}\varphi H)`$ (29)
The Wess-Zumino energy for the product of three matrices $`A,B,C`$ can be written as the sum of the energies for the actions for $`A,B`$ and $`C`$ plus another term:
$`S(ABC)=S(A)+S(B)+S(C){\displaystyle \frac{k}{4\mathrm{\Pi }}}{\displaystyle }d^2\xi Tr[\left(A^1_zA\right)\left(_{\overline{z}}B\right)B^1`$ (30)
$`+\left(B^1_{\overline{z}}B\right)\left(_zC\right)C^1+\left(A^1_{\overline{z}}A\right)B\left(_zC\right)C^1B^1]`$ (31)
Using our parametrization, the energy for the Wess-Zumino model takes the following local form:
$`_z_{\overline{z}}\varphi +2e^\varphi _zy_{\overline{z}}x=0`$ (32)
$`_z(_{\overline{z}}xe^\varphi )=_{\overline{z}}(_zye^\varphi )=0`$ (33)
The equations of motion for the Wess-Zumino model can be derived in these co-ordinates as
$`_z_{\overline{z}}\varphi +2e^\varphi _zy_{\overline{z}}x=0`$ (34)
$`_z(_{\overline{z}}xe^\varphi )=_{\overline{z}}(_zye^\varphi )=0`$ (35)
Consider special solutions
$$_{\overline{z}}x=\nu e^\varphi _zy=\mu e^\varphi $$
(36)
where $`\mu `$, $`\nu `$ are arbitrary constants. Then the system reduces to the Liouville system
$$_z_{\overline{z}}\varphi +2\mu \nu e^\varphi =0$$
(37)
## V Self-dual Chern Simons and Toda theory
In , Gerald Dunne defined the self-dual Chern-Simons equations over $`SU(N)`$ as:
$`_zA_{\overline{z}}_{\overline{z}}A_z+[A_{\overline{z}},A_z]`$ $`=`$ $`{\displaystyle \frac{2}{\kappa }}[\mathrm{\Psi }^{},\mathrm{\Psi }]`$ (38)
$`_{\overline{z}}\mathrm{\Psi }+[A_{\overline{z}},\mathrm{\Psi }]`$ $`=`$ $`0.`$ (39)
He showed that, with for a diagonal $`A`$ and upper triangular $`\mathrm{\Psi }`$,
$`A_z={\displaystyle A_i^\alpha H_\alpha }`$ (40)
$`\mathrm{\Psi }={\displaystyle \psi ^\alpha E_\alpha }`$ (41)
these equations combine to become those of the Toda model:
$`^2\varphi _\alpha ={\displaystyle \frac{2}{\kappa }}K_{\alpha \beta }\varphi _\beta .`$ (42)
where $`\mathrm{ln}\varphi _\alpha \left|\psi ^\alpha \right|^2`$. He also showed it is possible to make a gauge transformation $`u^1`$ which combines the self-dual Chern-Simons equations into a single equation:
$$_{\overline{z}}\chi =[\chi ^{},\chi ]$$
(43)
where
$$\chi =\sqrt{\frac{2}{\kappa }}u\mathrm{\Psi }u^1$$
(44)
Define
$`\stackrel{~}{A}_zA_z\sqrt{{\displaystyle \frac{2}{\kappa }}}\mathrm{\Psi }`$ (45)
$`\stackrel{~}{A}_{\overline{z}}A_{\overline{z}}+\sqrt{{\displaystyle \frac{2}{\kappa }}}\mathrm{\Psi }^{}`$ (46)
Then the self-dual Chern Simons equations imply that $`\stackrel{~}{A}`$ is flat. Trivialising $`\stackrel{~}{A}`$ as $`\stackrel{~}{A}=u^1du`$, and using the $`\chi `$ defined above, we find that
$`_zA_{\overline{z}}_{\overline{z}}A_z+[A_{\overline{z}},A_z]{\displaystyle \frac{2}{\kappa }}[\mathrm{\Psi }^{},\mathrm{\Psi }]=g^1\left(_{\overline{z}}\chi +_z\chi ^{}2[\chi ^{},\chi ]\right)g.`$ (47)
This shows that the equations (38) and (39) are equivalent to the single equation (43).
This equation may now be written as the harmonic map equation
$`_z\left(h^1_{\overline{z}}h\right)+_{\overline{z}}\left(h^1_zh\right)=0,`$ (48)
where $`hSU(N)`$ is related to $`\chi `$ as
$`h^1_zh=2\chi .`$ (49)
All SU(2) finite action harmonic maps have the form
$`h=h_0\left(2p1\right)`$ (50)
where $`p`$ is a holomorphic projection valued map
$`\left(1p\right)_zp=0.`$ (51)
With this condition, we can write a general $`p`$ in the defining representation for SU(2) as
$`p={\displaystyle \frac{MM^{}}{M^{}M}}`$ (52)
where
$`M=\left(\begin{array}{c}1\\ f(z)\end{array}\right)`$ (53)
for arbitrary $`f(z)`$. It can easily be checked that $`p`$ satisfies the correct projectivity, hermiticity and holomorphicity conditions.
Hence, we find that
$$p=\frac{1}{1+f\overline{f}}\left(\begin{array}{cc}1& \overline{f}\\ f& f\overline{f}\end{array}\right)$$
(54)
Explicitly, the $`\chi `$ become
$`\frac{1}{2}h^1_zh=\frac{1}{2}(2p1)2_zp=_zp`$ (55)
$`={\displaystyle \frac{f_z\overline{f}}{\left(1+f\overline{f}\right)^2}}\left(\begin{array}{cc}1& \frac{1}{f}\\ f& 1\end{array}\right).`$ (56)
The corresponding commutator is
$$[\chi ,\chi ^{}]=\frac{_z\overline{f}_{\overline{z}}f}{\left(1+f\overline{f}\right)^3}\left(\begin{array}{cc}1f\overline{f}& 2\overline{f}\\ 2f& 1+f\overline{f}\end{array}\right)$$
(57)
Dunne showed that this can be diagonalised by $`uSU(2)`$
$$u=\frac{1}{\sqrt{1+f\overline{f}}}\left(\begin{array}{cc}1& \overline{f}\\ f& 1\end{array}\right),$$
(58)
so that
$$u^1[\chi ,\chi ^{}]u=\frac{_z\overline{f}_{\overline{z}}f}{\left(1+f\overline{f}\right)^2}\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right).$$
(59)
Using (44), we see that $`\mathrm{ln}\varphi _\alpha \left|\psi ^\alpha \right|^2`$ is the quantity
$$\frac{_z\overline{f}_{\overline{z}}f}{\left(1+f\overline{f}\right)^2}$$
where $`\varphi `$ is the Toda field. It above can also be written as
$$_z_{\overline{z}}\mathrm{ln}(1+f\overline{f})=_z_{\overline{z}}\mathrm{ln}det(M^{}M)$$
(60)
This is the form of the general solution of the classical Liouville equation .
## VI Harmonic Maps and Toda Systems
We seek a broader understanding of the relationship of harmonic maps to other integrable systems. In the process we construct a general scheme for reducing certain Harmonic maps into Lie groups to the Toda system and examine the relationships with other models discussed in the previous chapter.
We look mainly at harmonic maps into SL(2,R) or equivalently, SU(1,1), a choice which allows us to compare our constructions to both the work of Dunne and O’Raifertaigh et. al. . Using analysis described by Uhlenbeck in for harmonic maps into unitary groups, we find that the work of Dunne can be used to find solutions for the Liouville system. The second section of this chapter employs the field-theoretic approach of O’Raifertaigh et. al. using a decomposition of the group SU(1,1) to demonstrate an alternative reduction to the Toda model.
We then construct a new framework which puts these reductions in a single context. A requirement that the fields in the Lax pair satisfy certain requirements in their Lie algebraic structure is found to be behind the conditions and Ansatze used in , and , and others.
There exists a relatively straightforward way to associate a harmonic map and the corresponding harmonic map (chiral model) equations with a CAT theory. This method also holds for reducing the non-affine algebra (and the corresponding system of equations) to the simple Toda theory.
The correspondence relies on the following proposition, first proved in \[U2\]:
Proposition 1 If $`A=A_zdz+A_{\overline{z}}d\overline{z}`$ and $`B=B_zdz+B_{\overline{z}}d\overline{z}`$ satisfy
$$[_{\overline{z}}+A_{\overline{z}}+\lambda B_{\overline{z}},_z+A_z+\lambda ^1B_z]=0,$$
(61)
then there exists $`u`$ with
$$\stackrel{~}{A}=u^1(Au+du),\stackrel{~}{B}=u^1Bu$$
(62)
and
$$[_{\overline{z}}+(1\lambda )\stackrel{~}{A}_{\overline{z}},_z+(1\lambda ^1)\stackrel{~}{A}_z]=0.$$
(63)
Furthermore, $`2\stackrel{~}{A}=s^1ds`$ where $`s`$ is harmonic.
Proof: The necessary transformations correspond to trivialising (61) at $`\lambda =1`$, with the gauge transformation $`u`$. From , since $`𝐂`$ is simply connected, we find that $`s`$ is a harmonic map where $`\stackrel{~}{A}=s^1ds/2`$.
Notice that (20) can be written in the correct format to allow use of Proposition 1, where
$`A_z`$ $`=_z\mathrm{\Phi }+e^\mathrm{\Phi }E_+e^\mathrm{\Phi }`$ (64)
$`A_z`$ $`=_{\overline{z}}\mathrm{\Phi }+e^\mathrm{\Phi }E_{}e^\mathrm{\Phi }`$ (65)
$`B_z`$ $`=e^\mathrm{\Phi }E_{}e^{Phi}`$ (66)
$`B_{\overline{z}}`$ $`=e^\mathrm{\Phi }E_+e^\mathrm{\Phi }.`$ (67)
We write the Lax pair of our harmonic map as
$$(_{\overline{z}}(1\lambda )u^1e^\mathrm{\Phi }E_+e^\mathrm{\Phi }u,_z(1\lambda ^1)u^1e^\mathrm{\Phi }E_{}e^\mathrm{\Phi }u),$$
(68)
an equation which holds for all order of $`\lambda `$.
## VII A condition for the equivalence of the harmonic map and Conformal Affine Toda equations
Not all harmonic maps are Toda maps. We now show that Toda maps, and more generally, CAT maps, are a subset of harmonic maps into the appropriate Lie group.
As a brief introduction on the way to a more general result for all semisimple Lie groups, we now describe a condition under which a harmonic map into SL(2,C)’ (the subgroup of constant loops in the centrally extended loop group) gives rise to the CAT system of equations.
Theorem 2 If $`s`$ is a harmonic map into SL(2,C), then the harmonic map equations for $`s`$ are gauge equivalent to the Conformal Affine Toda equations (92) for $`\varphi ,\eta ,\xi `$ if the fields $`s^1_zs,s^1_{\overline{z}}s`$ can be simultaneously diagonalised into lower and upper triangular matrices respectively.
Proof: Let $`\stackrel{~}{A}_z=s^1_zs/2,\stackrel{~}{A}_{\overline{z}}=s^1_{\overline{z}}s/2`$. Thus we can write, for some $`\theta \widehat{sl_2}^{}`$,
$$\stackrel{~}{A}_z=e^\theta fE_+e^\theta ,\stackrel{~}{A}_z=e^\theta f^{}E_{}e^\theta $$
(69)
By proposition 1, the harmonic map equations for $`s`$ are described by the vanishing of the Lax pair in (63) for all $`\lambda C^{}`$. In terms of the connection given above, this provides the following identites as coefficients for $`\lambda ^1,\lambda `$ respectively:
$`_{\overline{z}}fE_+f[_{\overline{z}}\theta ,E_+]2f^{}fH=0`$ (70)
$`_zf^{}E_++f^{}[_z\theta ,E_{}]2f^{}fH=0`$ (71)
Writing $`_z\theta ,_{\overline{z}}\theta `$ in terms of its algebraic decomposition
$`_z\theta =\alpha E_++\beta H+\gamma E_{}+\delta c`$ (72)
$`_{\overline{z}}\theta =\alpha ^{}E_++\beta ^{}H+\gamma ^{}E_{}+\delta ^{}c`$ (73)
we see that the following relationships are specified:
$`_zf^{}+\beta f^{}=0\alpha ff^{}f=0`$ (74)
$`_{\overline{z}}f\beta ^{}f=0\gamma ^{}f^{}f^{}f=0`$ (75)
Using this information, we write
$`_z\theta =f^{}E_+_z\left(lnf^{}\right)H+\gamma E_{}+\delta c`$ (76)
$`_{\overline{z}}\theta =fE_{}+_{\overline{z}}\left(lnf\right)H+\alpha ^{}E_++\delta ^{}c.`$ (77)
Apply Proposition 1 in exactly the reverse manner using the inverse of the gauge transformation given by the trivialisation at the identity. In other words, transform to
$`A=u\stackrel{~}{A}u^1duu^1`$ (78)
$`B=u\stackrel{~}{B}u=u\stackrel{~}{A}u^1`$ (79)
where here we have $`u=e^\theta `$. Define the quantities
$`\varphi _1`$ $`=lnf;\eta _1`$ $`=ln\gamma +\varphi _1`$ (80)
$`\varphi _2`$ $`=lnf^{};\eta _2`$ $`=ln\alpha ^{}+\varphi _2.`$ (81)
We find the connection becomes
$`(_{\overline{z}}+A_{\overline{z}}+\lambda B_{\overline{z}},_z+A_z+\lambda ^1B_z)`$ (82)
$`=(_{\overline{z}}_{\overline{z}}\varphi _1He^{\eta _1\varphi _1}E_+\delta ^{}c\lambda e^{\varphi _1}E_{},_z+_z\varphi _2He^{\eta _2\varphi _2}E_{}\delta c\lambda ^1e^{\varphi _2}E_+).`$ (83)
The curvature is given by
$`\left(_z_{\overline{z}}\left(\varphi _1+\varphi _2\right)+2e^{\varphi _1+\varphi _2}2e^{\eta _1+\eta _2\varphi _1\varphi _2}\right)H`$ (84)
$`+\left(_ze^{\eta _1\varphi _1}+\lambda ^1_{\overline{z}}\varphi _1e^{\varphi _2}\lambda ^1_{\overline{z}}\varphi _2e^{\varphi _2}_z\varphi _2e^{\eta _1\varphi _1}\right)E_+`$ (85)
$`+\left(_{\overline{z}}e^{\eta _2\varphi _2}+\lambda _z\varphi _2e^{\varphi _1}\lambda ^1_z\varphi _1e^{\varphi _1}+_{\overline{z}}\varphi _1e^{\eta _2\varphi _2}\right)E_{}`$ (86)
$`+\left(_z\delta ^{}+_{\overline{z}}\delta +e^{\varphi _1+\varphi _2}\right)c.`$ (87)
We now introduce the variables $`\xi _1,\xi _2`$ where $`\delta ^{}=_{\overline{z}}\xi _1,\delta _2=_z\xi _2`$. We find that, by writing
$`\eta =\eta _1+\eta _2;\varphi =\varphi _1+\varphi _2;`$ (88)
$`\xi =\xi _1+\xi _2+\varphi `$ (89)
this system is equivalent to the CAT set:
$`_z_{\overline{z}}\varphi `$ $`=`$ $`e^{2\varphi }e^{2\eta 2\varphi }`$ (90)
$`_z_{\overline{z}}\eta `$ $`=`$ $`0`$ (91)
$`_z_{\overline{z}}\xi `$ $`=`$ $`e^{2\eta 2\varphi }.`$ (92)
Note that a reduction of our result to the unextended sl(2) algebra will lead to a subset of the CAT equations, namely equations (90) and (91). Although equation (92) seems somewhat superfluous to requirements, in that it can be safely omitted while preserving the conformal invariance, the central extension plays an important role by including the spectral parameter into the algebra. Also, the conformal invariance of the CAT theory can be described by the invariance of the harmonic map equations for the system described above under the transformation $`f\stackrel{~}{f}=g(z)f`$ .
## VIII Algebraic Reduction of Harmonic Maps to Toda Systems
In the last section we reduced the harmonic map equations to the Conformal Affine Toda systems for SL(2). In this section, we will investigate harmonic maps into SU(1,1) and see how we can identify these with solutions of the Toda equations. Later we will generalise this to include all semi-simple Lie groups.
Let us first construct the analog of Uhlenbeck’s uniton for the case of a harmonic map into SU(1,1). We will begin by generalising the uniton construction of to the indefinite $`U(q,Nq)`$ form. Again, we consider maps into the Grassmannian of $`k`$ planes which is a geodesic submanifold.
Definition: An $`n,q`$-uniton is a harmonic map $`s:\mathrm{\Omega }U(q,Nq)`$ which has an extended solution
$$E_\lambda :𝐂^{}\times \mathrm{\Omega }G=GL(N,C)$$
(93)
with
$`(a)`$ $`E_\lambda ={\displaystyle _{\alpha =0}^n}T_\alpha \lambda ^\alpha forT_\alpha :\mathrm{\Omega }gl(N,C)`$ (94)
$`(b)`$ $`E_1=I`$ (95)
$`(c)`$ $`E_1=Qs^1forQSU(q,Nq)constant`$ (96)
$`(d)`$ $`(E_{\overline{\lambda }})^{}J=J(E_{\lambda ^1})^1`$ (97)
and $`J`$ is the automorphism which defines the real form for $`U(q,Nq)`$
Proposition 3: $`s:\mathrm{\Omega }SU(q,Nq)`$ is a 1-q-uniton (a holomorphic map) if $`s=Q(2p1)`$ for constant $`QSU(q,Nq)`$, where $`p`$ satisfies $`p^{}J=Jp`$, $`p^2=p`$, and $`(1p)_{\overline{z}}p=0`$
Proof: With $`E_\lambda =T_0+\lambda T_1`$ and $`E_1=T_0+T_1=I`$, we have $`E_\lambda =T_0+\lambda (IT_0)`$.
The reality condition (d) above tells us that
$`(IT_0^{})JT_0`$ $`=`$ $`0`$ (98)
$`T_0^{}JT_0+(IT_0^{})J(IT_0)`$ $`=`$ $`J`$ (99)
Combining these equations gives $`T_0^{}J=JT_0`$ and $`T_0^2=T_0`$. We can identify $`T_0`$ with $`p`$. As in the previous chapter the necessary condition for $`E_\lambda `$ to be an extended solution for a harmonic map as is that $`(1p)_{\overline{z}}p=0`$
With this condition, we can write a general $`p`$ in the defining representation for SU(1,1) as
$$p=M(M^{}JM)^1M^{}J$$
(100)
where, as before,
$$M=\left(\begin{array}{c}1\\ f(x_{})\end{array}\right)$$
(101)
for arbitrary $`f(x_{})`$.
Hence, we find that, in this case
$$p=\frac{1}{1f\overline{f}}\left(\begin{array}{cc}1& \overline{f}\\ f& f\overline{f}\end{array}\right)$$
(102)
The simplest solutions of the harmonic map equation are now given by the holomorphic maps:
$`\frac{1}{2}g^1_zg=\frac{1}{2}(2p1)2_zp=_zp`$ (103)
$`={\displaystyle \frac{f_z\overline{f}}{\left(1f\overline{f}\right)^2}}\left(\begin{array}{cc}1& \frac{1}{f}\\ f& 1\end{array}\right)`$ (104)
We can now find a $`uSU(1,1)`$
$$u=\frac{1}{\sqrt{1f\overline{f}}}\left(\begin{array}{cc}1& \overline{f}\\ f& 1\end{array}\right)$$
(105)
such that
$$u^1[A_+,A_{}]u=\frac{_z\overline{f}_{\overline{z}}f}{\left(1f\overline{f}\right)^2}\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)$$
(106)
This diagonalisation is a critical step, and as was shown in the previous chapter for SU(2), the magnitude of this quantity should be a solution of the appropriate Toda equation. It can be rewritten
$$u=\frac{1}{\sqrt{1f\overline{f}}}\left(\begin{array}{cc}1& \overline{f}\\ f& 1\end{array}\right)$$
(107)
and we find that the magnitude of this expression for the diagonalised commutator above is just
$$\mathrm{ln}\varphi =_z_{\overline{z}}\mathrm{ln}(1f\overline{f})=_z_{\overline{z}}\mathrm{ln}det(M^{}JM).$$
(108)
## IX Harmonic Maps and Liouville Theory
Let us now examine the reduction of the SU(1,1) harmonic map to the Toda system using a decomposition into local fields.
Recall that the energy or action for $`g:\mathrm{\Omega }SU(1,1)`$ is
$$S(g)=\frac{k}{8\pi }d^2zTr\left[(g^1_zg)(g^1_{\overline{z}}g)\right]$$
(109)
As in the previous chapter, where we saw a similar reduction to analyse Wess-Zumino-Witten theory, we use a Gauss decomposition of $`SU(1,1)`$. For regular $`g`$ in a neighbourhood of the identity, we write:
$$g=ABC$$
(110)
where
$`A=\left(\begin{array}{cc}1& ix\\ 0& 1\end{array}\right)=\mathrm{exp}(ixE_+)`$ (111)
$`C=\left(\begin{array}{cc}1& 0\\ iy& 1\end{array}\right)=\mathrm{exp}(iyE_{})`$ (112)
$`B=\left(\begin{array}{cc}\mathrm{exp}(\frac{1}{2}\varphi )& 0\\ 0& \mathrm{exp}(\frac{1}{2}\varphi )\end{array}\right)=\mathrm{exp}(\frac{1}{2}\varphi H)`$ (113)
With this parametrization, the energy for the harmonic map can be written in terms of three real fields as
$$S(g)=S(x,y,\varphi )=\frac{k}{8\pi }d^2z\left[\frac{1}{2}_z\varphi _{\overline{z}}\varphi e^\varphi (_zx_{\overline{z}}y+_zy_{\overline{z}}x)\right]$$
(114)
A local form of the equations of motion for the harmonic map can now be derived from this:
$`_z_{\overline{z}}\varphi e^\varphi (_zx_{\overline{z}}y+_zy_{\overline{z}}x)=0`$ (115)
$`_z(_{\overline{z}}xe^\varphi )+_{\overline{z}}(_zxe^\varphi )=0`$ (116)
$`_z(_{\overline{z}}ye^\varphi )+_{\overline{z}}(_zye^\varphi )=0.`$ (117)
We consider the solution for the above:
$$_zx=f(z)e^\varphi ,_zy=g(z)e^\varphi .$$
(118)
Then the system reduces to the following form of the Liouville system:
$`_z_{\overline{z}}\varphi +M(z,\overline{z})e^\varphi =0`$ (119)
$`f\left(z\right)\overline{g}\left(\overline{z}\right)\overline{f}\left(\overline{z}\right)g\left(z\right)M=0.`$ (120)
## X A useful result
We have seen the reduction of the harmonic maps into Lie groups to the Toda model carried out in a number of different ways, and there are others, such as , which we have not alluded to. We now ask ourselves if there is any underlying link in the manner in which these reductions are carried out.
We find that these reductions can be observed to be different strategies for taking advantage of the following general result:
Result 4 Any pencil of connections $`(A_z+\lambda A_z^{},A_{\overline{z}}+\lambda ^1A_{\overline{z}}^{})`$ gives rise to the Toda system of equations for the corresponding semi-simple Lie group if the $`A_z,A_{\overline{z}}`$ are upper or lower triangular respectively and the $`A_z^{},A_{\overline{z}}^{}`$ are off-diagonal.
Proof:
We will begin with the most general case and reduce using our stated condition as needed. We see that most of the work is accomplished by the algebraic structure and the parameter $`\lambda `$.
We begin with the stated condition that
$$[_z+A_z+\lambda A_z^{},_{\overline{z}}+A_{\overline{z}}+\lambda ^1A_{\overline{z}}^{}]=0$$
(121)
where $`A`$ and $`A^{}`$ take values in the Chevalley basis of the Lie algebra, i.e.
$$\begin{array}{c}A_z=\underset{\alpha \mathrm{\Phi }^+}{}g_+^\alpha E_\alpha +\underset{\alpha \mathrm{\Phi }^+}{}f_+^\alpha E_\alpha +\underset{\alpha \mathrm{\Phi }^+}{}h_+^\alpha H_\alpha \\ A_{\overline{z}}=\underset{\alpha \mathrm{\Phi }^+}{}g_{}^\alpha E_\alpha +\underset{\alpha \mathrm{\Phi }^+}{}f_{}^\alpha E_\alpha +\underset{\alpha \mathrm{\Phi }^+}{}h_{}^\alpha H_\alpha \end{array}$$
(122)
and
$$\begin{array}{c}A_z^{}=\underset{\alpha \mathrm{\Phi }^+}{}g_+^\alpha E_\alpha +\underset{\alpha \mathrm{\Phi }^+}{}f_+^\alpha E_\alpha +\underset{\alpha \mathrm{\Phi }^+}{}h_+^\alpha H_\alpha \\ A_{\overline{z}}^{}=\underset{\alpha \mathrm{\Phi }^+}{}g_{}^\alpha E_\alpha +\underset{\alpha \mathrm{\Phi }^+}{}f_{}^\alpha E_\alpha +\underset{\alpha \mathrm{\Phi }^+}{}h_{}^\alpha H_\alpha \end{array}$$
(123)
Here $`\mathrm{\Phi }^+`$ denotes the set of positive roots, and $`H_\alpha ,E_{\pm \alpha }`$ are the Cartan subalgebra and step generators of the Chevalley basis respectively.
In terms of the first power of of $`\lambda `$ and the basis of the Lie algebra, the set of equations reduces to:
$`_{\overline{z}}g_+^\alpha {\displaystyle \underset{\beta }{}}K_{\beta \alpha }h_{}^\beta g_+^\alpha +{\displaystyle \underset{\beta }{}}K_{\beta \alpha }h_+^\beta g_{}^\alpha =0`$ (124)
$`_{\overline{z}}f_+^\alpha +{\displaystyle \underset{\beta }{}}K_{\beta \alpha }h_{}^\beta f_+^\alpha {\displaystyle \underset{\beta }{}}K_{\beta \alpha }h_+^\beta f_{}^\alpha =0`$ (125)
$`_{}h_+^\alpha +g_+^\alpha f_{}^\alpha f_+^\alpha g_{}^\alpha =0`$ (126)
and their dual, where $`K`$ is the classical Cartan matrix for the Lie algebra. In the case where $`A^{}`$ is off-diagonal, all the $`h^{}`$ vanish, so we find:
$`_zh_{}^\alpha _{\overline{z}}h_+^\alpha +f_{}^\alpha g_+^\alpha f_+^\alpha g_{}^\alpha +f_{}^\alpha g_+^\alpha f_+^\alpha g_{}^\alpha =0`$ (127)
The part of (121) independent of $`\lambda `$ contributes the following (and its dual:
$$_zg_{}^\alpha _{\overline{z}}g_+^\alpha \underset{\beta }{}K_{\beta \alpha }h_{}^\beta g_+^\alpha +\underset{\beta }{}K_{\beta \alpha }h_+^\beta g_{}^\alpha \underset{\beta }{}K_{\beta \alpha }h_{}^\beta g_+^\alpha +\underset{\beta }{}K_{\beta \alpha }h_+^\beta g_{}^\alpha =0$$
(128)
as well as
$$_zh_{}^\alpha _{\overline{z}}h_+^\alpha +f_{}^\alpha g_+^\alpha f_+^\alpha g_{}^\alpha +f_{}^\alpha g_+^\alpha f_+^\alpha g_{}^\alpha =0$$
(129)
When $`A`$ is triangular, $`g_+^\alpha =0`$ and when $`A^{}`$ is off-diagonal, $`h^{}=0`$ and we find $`g_{}^\alpha =0`$. Equation (126) and its dual now give us:
$`_z\left(\mathrm{ln}g_{}^\alpha \right)={\displaystyle \underset{\beta }{}}K_{\beta \alpha }h_+^\beta `$ (130)
$`_z\left(\mathrm{ln}f_+^\alpha \right)={\displaystyle \underset{\beta }{}}K_{\beta \alpha }h_{}^\beta `$ (131)
Combining these equations we find
$$_z_{\overline{z}}\left[\mathrm{ln}(g_+^\alpha f_{}^\alpha )\mathrm{ln}\left(g_{}^\alpha f_+^\alpha \right)\right]=2\underset{\beta }{}K_{\beta \alpha }\left(g_+^\alpha f_{}^\alpha g_{}^\alpha f_+^\alpha \right).$$
(132)
With the further result that
$$_{\overline{z}}\mathrm{ln}\left(g_+^\alpha f_+^\alpha \right)=_z\mathrm{ln}\left(f_{}^\alpha g_{}^\alpha \right)=0,$$
(133)
we recognise (132) as the affine toda equations:
$$_z_{\overline{z}}\varphi _\alpha +\underset{\beta }{}K_{\beta \alpha }\left(\eta _+^\alpha e^{\varphi _\alpha }\eta _{}^\alpha e^{\varphi _\alpha }\right).$$
(134)
Here
$$\varphi _\alpha =\mathrm{ln}\left(\frac{f_{}^\alpha }{f_+^\alpha }\right)$$
(135)
and $`\eta _+^\alpha `$ , $`\eta _{}^\alpha `$ are arbitrary holomorphic and anti-holomorphic functions respectively.
## XI Some conclusions
All of the reductions to Toda systems we have seen use the above Result 4, although this is far from clear from a first examination. An understanding of why this follows from the fact that for a large class of harmonic maps, the paramatrised Lax pair can be gauge transformed into our required form. We need the following result.
Result 5: If the commutator $`[A_z,A_{\overline{z}}]`$ is diagonal, then either $`A_z,A_{\overline{z}}`$ are off-diagonal or the commutator is zero
Proof: Using the expansion of $`A`$ as in (122), we find that since the off-diagonal elements of the commutator are given for all $`\alpha ,\beta `$ by
$`h_+^\beta g_{}^\alpha h_{}^\beta g_+^\alpha =0`$ (136)
$`h_+^\beta f_{}^\alpha h_{}^\beta f_+^\alpha =0`$ (137)
We find that the only case when all $`h^\beta `$ do not vanish is the trivial case.
Now when we gauge transform by a group element $`u`$ which diagonalises the hermitian quantity $`[A_z,A_{\overline{z}}]`$ with $`u^1du`$ triangular, our connection is in the appropriate form to apply the Result 4. In the work on self-dual Chern-Simons theory by Dunne , it can be seen that for the SU(N) uniton solutions, the necessary gauge transformation is in the required form. These uniton solutions are therefore Toda systems. Moreover, Guest , selected a priori a connection equivalent to our form for his reduction of the chiral model to the Toda lattice.
O’Raifeartaigh et al carry out essentially this same diagonalisation process but in a different setting. To see this, use the decomposition in (110) and write the commutator of the connection of Section V of this chapter as
$$C^1B^1[\stackrel{~}{A}_+,\stackrel{~}{A}_{}]BC$$
(138)
where
$$\stackrel{~}{A}=A^1dA+(dB)B^1+B(dC)C^1B^1$$
(139)
Now we see that the quantity $`BC`$ corresponds to the $`u`$ in (107) which diagonalises $`[\stackrel{~}{A}_+,\stackrel{~}{A}_{}]`$. In light of the conditions (118), the diagonal component
$$\left(_zx_{\overline{z}}y_zx_{\overline{z}}y\right)e^\varphi $$
(140)
is just the Liouville field $`e^\varphi `$ up to a holomorphic function. This corresponds to the SU(1,1) Toda solution result in (108) and confirms the solution found by Fujii . |
warning/0506/hep-ex0506036.html | ar5iv | text | # Determination of |𝑉_{𝑢𝑏}| from Measurements of the Electron and Neutrino Momenta in Inclusive Semileptonic 𝐵 Decays
## I Acknowledgments
We are grateful for the excellent luminosity and machine conditions provided by our PEP-II colleagues, and for the substantial dedicated effort from the computing organizations that support BABAR. The collaborating institutions wish to thank SLAC for its support and kind hospitality. This work is supported by DOE and NSF (USA), NSERC (Canada), IHEP (China), CEA and CNRS-IN2P3 (France), BMBF and DFG (Germany), INFN (Italy), FOM (The Netherlands), NFR (Norway), MIST (Russia), and PPARC (United Kingdom). Individuals have received support from CONACyT (Mexico), Marie Curie EIF (European Union), the A. P. Sloan Foundation, the Research Corporation, and the Alexander von Humboldt Foundation. Finally, we would like to thank the many theorists with whom we have had valuable discussions, and further thank M. Neubert, B. Lange and G. Paz for making available for our use a computer code implementing their calculations. |
warning/0506/hep-th0506077.html | ar5iv | text | # Contents
## 1 Introduction and summary
In the last few years the instabilities associated with open string tachyons have been studied extensively and have become reasonably well understood . The instabilities associated with closed string tachyons have proven to be harder to understand. For the case of localized closed string tachyons – tachyons that live on subspaces of spacetime – there are now plausible conjectures for the associated instabilities and a fair amount of circumstantial evidence for them .
The bulk tachyon of the closed bosonic string is the oldest known closed string tachyon. It remains the most mysterious one and there is no convincing analysis of the associated instability. The analogy with open strings, however, suggests a fairly dramatic possibility. In open bosonic string in the background of a spacefilling D-brane, the tachyon potential has a critical point that represents spacetime without the D-brane and thus without physical open string excitations. In an analogous closed string tachyon vacuum one would expect no closed string excitations. Without gravity excitations spacetime ceases to be dynamical and it would seem that, for all intents and purposes, it has dissappeared.
There has been no consensus that such a closed string tachyon vacuum exists. In fact, no analysis of the closed string tachyon potential (either in the CFT approach or in the SFT approach) has provided concrete evidence of a vacuum with non-dynamical spacetime. Since the analogous open string tachyon vacuum shows up quite clearly in the open string field theory computation of the potential it is natural to consider the corresponding calculation in closed string field theory (CSFT) .
The quadratic and cubic terms in the closed string tachyon potential are well known :
$$\kappa ^2𝕍_0^{(3)}=t^2+\frac{6561}{4096}t^3,(\alpha ^{}=2).$$
(1.1)
These terms define a critical point analogous to the one that turns out to represent the tachyon vacuum in the open string field theory. In open string field theory higher level computations make the vacuum about 46% deeper. Since CSFT is nonpolynomial, it is natural to investigate the effect of the quartic term in the potential. This term was found to be
$$\kappa ^2V_0^{(4)}=3.0172t^4.$$
(1.2)
This term is so large and negative that $`𝕍_0^{(3)}+V_0^{(4)}`$ has no critical point. In fact, the quartic term in the effective tachyon potential (obtained by integrating out massive fields) is even a bit larger . The hopes of identifying a reliable critical point in the closed string tachyon potential were dashed<sup>1</sup><sup>1</sup>1In the effective open string tachyon potential a negative quartic term also destroys the cubic critical point. Nevertheless, the critical point can be gleaned using Pade-approximants . For closed strings, however, the quartic term is too large: for a potential $`v(t)=v_2t^2+v_3t^3+v_4t^4`$, with $`v_2,v_4<0`$, the approximant formed by the ratio of a cubic and a linear polynomial fails to give a critical point when $`v_2v_4v_3^2`$..
Recent developments inform our present analysis. The tachyon potential must include all fields that are sourced by the zero-momentum tachyon. As discussed in , this includes massless closed string states that are built from ghost oscillators, in particular, the zero-momentum ghost-dilaton state $`(c_1c_1\overline{c}_1\overline{c}_1)|0`$. The search for a critical point cannot be carried out consistently without including the ghost dilaton. Computations of quartic vertices coupling dilatons, tachyons, and other massive fields are now possible due to the work of Moeller and have been done to test the marginality of matter and dilaton operators .
As we explain now, ghost-dilaton couplings to the tachyon restore the critical point in the potential. The key effect can be understood from the cubic and quartic couplings
$$\kappa ^2V(t,d)=\frac{27}{32}td^2+\mathrm{\hspace{0.17em}3.8721}t^3d+\mathrm{}.$$
(1.3)
The cubic coupling plays no role as long as we only consider cubic interactions: $`d`$ can be set consistently to zero. The quartic coupling is linear in $`d`$. Once included, the equation of motion for the dilaton can only be satisfied if the dilaton acquires an expectation value. Solving for the dilaton one finds $`d=2.2944t^2`$ and substituting back,
$$\kappa ^2V(t,d)=4.4422t^5+\mathrm{}$$
(1.4)
This positive quintic term suffices to compensate the effects of (1.2) and restores the critical point. Our computations include additional couplings and the effect of massive fields as well. The critical point persists and may be reliable, although more work is needed to establish this convincingly.
In order to interpret the critical point we raise and answer a pair of questions. The ghost-dilaton has a positive expectation value at the critical point. Does this correspond to stronger or weaker string coupling ? We do a detailed comparison of quadratic and cubic terms in the closed string field theory action and in the low-energy effective field theory action. The conclusion is that the positive dilaton expectation value corresponds to stronger coupling. In our solution the ghost-dilaton is excited but the scalar operator $`c\overline{c}X\overline{}X`$, sometimes included in the dilaton vertex operator, is not. We ask: Is the string metric excited? Is the Einstein metric excited? These questions are only well-defined at the linearized level, but the answers are clear: the string metric does not change, but the Einstein metric does. We take the opportunity to explain the relations between the four kinds of “dilatons” that are used in the literature: the ghost-dilaton, the matter-dilaton, the dilaton, and the dilaton of the older literature. It is noted that one cannot define unambiguously a dilaton vertex operator unless one specifies which metric is left invariant; conversely, the metric vertex operator is only determined once one specifies which dilaton is left invariant.
In a companion paper we attempted to gain insight into the tachyon vacuum by considering the rolling solutions<sup>2</sup><sup>2</sup>2Rolling solutions have long been considered using Liouville field theory to provide conformal invariant sigma model with spacetime background fields that typically include a linear dilaton and a constant string metric . of a low-energy effective action for the string metric $`g_{\mu \nu }`$, the tachyon $`T`$, and the dilaton $`\mathrm{\Phi }`$:
$$S_\sigma =\frac{1}{2\kappa ^2}d^Dx\sqrt{g}e^{2\mathrm{\Phi }}\left(R+4(_\mu \mathrm{\Phi })^2(_\mu T)^22V(T)\right).$$
(1.5)
This action, suggested by the beta functions of sigma models with background fields , is expected to capture at least some of the features of string theory solutions. The potential is tachyonic: $`V(T)=\frac{1}{2}m^2T^2+𝒪(T^3)`$, but is otherwise left undetermined. We found that solutions in which the tachyon begins the rolling process always have constant string metric for all times – consistent with the type of the SFT critical point. The dilaton, moreover, grows in time throughout the evolution – consistent with the larger dilaton vev in the SFT critical point. Rather generally, the solution becomes singular in finite time: the dilaton runs to infinity and the string coupling becomes infinite. Alternatively, the Einstein metric crunches up and familiar spacetime no longer exists. This seems roughly consistent with the idea that the tachyon vacuum does not have a fluctuating spacetime.
Perhaps the most subtle point concerns the value of the on-shell action. In the open string field theory computation of the tachyon potential, the value of the action (per unit spacetime volume) is energy density. The tachyon conjectures are in fact formulated in terms of energy densities at the perturbative and the non-perturbative vacuum . Since the tree-level cosmological constant in closed string theory is zero, the value of the action at the perturbative closed string vacuum is zero. We ask: What is the value of the potential, or action (per unit volume) at the critical point ? The low-energy action (1.5) suggests a surprising answer. Consider the associated equations of motion:
$$\begin{array}{cc}\hfill R_{\mu \nu }+2_\mu _\nu \mathrm{\Phi }(_\mu T)(_\nu T)& =0,\hfill \\ \hfill ^2T2(_\mu \mathrm{\Phi })(^\mu T)V^{}(T)& =0,\hfill \\ \hfill ^2\mathrm{\Phi }2(_\mu \mathrm{\Phi })^2V(T)& =0.\hfill \end{array}$$
(1.6)
If the fields acquire constant expectation values we can satisfy the tachyon equation if the expectation value $`T_{}`$ is a critical point of the potential: $`V^{}(T_{})=0`$. The dilaton equation imposes an additional constraint: $`V(T_{})=0`$, the potential must itself vanish. This is a reliable constraint that follows from a simple fact: in the action the dilaton appears without derivatives only as a multiplicative factor. This fact remains true after addition of $`\alpha ^{}`$ corrections of all orders. It may be that $`V(T)`$ has a critical point $`T_0`$ with $`V(T_0)<0`$, but this cannot be the tachyon vacuum. The effective field equations imply that a vacuum with spacetime independent expectation values has zero action.
The action (1.5) can be evaluated on-shell using the equations of motion. One finds
$$S_{onshell}=\frac{1}{2\kappa ^2}d^{d+1}x\sqrt{g}e^{2\mathrm{\Phi }}\left(4V(T)\right).$$
(1.7)
In rolling solutions the action density changes in time but, as $`\mathrm{\Phi }\mathrm{}`$ at late times the action density goes to zero . This also suggests that the tachyon vacuum is a critical point with zero action.
In Figure 1 we present the likely features of the tachyon potential. The unstable perturbative vacuum $`T=0`$ has zero cosmological constant, and so does the tachyon vacuum $`T=\mathrm{}`$. The infinite value of $`T`$ is suggested by the analogous result in the effective open string theory tachyon potential (see conclusions). In SFT the tachyon vacuum appears for finite values of the fields, but the qualitative features would persist. The potential is qualitatively in the class used in cyclic universe models .
In our calculations we find some evidence that the action density, which is negative, may go to zero as we increase the accuracy of the calculation. To begin with, the value $`\mathrm{\Lambda }_0`$ of the action density at the critical point of the cubic tachyon potential (1.1) may be argued to be rather small. It is a cosmological term about seventy times smaller than the “canonical” one associated with $`D=2`$ non-critical string theory (see , footnote 5). Alternatively, $`\mathrm{\Lambda }_0`$ is only about 4% of the value that would be obtained using the on-shell coupling of three tachyons to calculate the cubic term. The inclusion of cubic interactions of massive fields makes the action density about 10% more negative. This shift, smaller than the corresponding one in open string field theory, is reversed once we include the dilaton quartic terms. In the most accurate computation we have done, the action density is down to 60% of $`\mathrm{\Lambda }_0`$. Additional computations are clearly in order.
As a by-product of our work, we investigate large dilaton deformations in CSFT. For ordinary marginal deformations the description reaches an obstruction for some finite critical value of the string field marginal parameter . The critical value is stable under level expansion, and the potential for the marginal field (which should vanish for infinite level) is small. For the dilaton, however, the lowest-order obstruction is not present . We carry this analysis to higher order and no reliable obstructions are found: critical values of the dilaton jump wildly with level and appear where the dilaton potential is large and cannot be trusted. This result strengthens the evidence that CSFT can describe backgrounds with arbitrarily large variations in the string coupling. If the infinite string coupling limit is also contained in the configuration space it may be possible to define M-theory using type IIA superstring field theory.
Let us briefly describe the contents of this paper. In section 2 we reconsider the universality arguments that require the inclusion of the ghost-dilaton, exhibit a world-sheet parity symmetry that allows a sizable truncation of the universal space, and note that universality may apply in circumstances significantly more general that originally envisioned . Our computational strategy for the tachyon potential, motivated by the results of , goes as follows. We compute all quadratic and cubic terms in the potential including fields up to level four. We then begin the inclusion of quartic terms and obtain complete results up to quartic interactions of total level four. The results make it plausible that a critical point exists and that the value of the action density decreases in magnitude as the accuracy improves. In section 3 we find the linearized relations between the metric, dilaton, and tachyon closed string fields and the corresponding fields in the sigma-model approach to string theory. These relations allow us to establish that the dilaton vev at the critical point represents an increased string coupling and that the string field at the critical point does not have a component along the vertex operator for the string metric. We discuss the vertex operators associated with the various definitions of the dilaton, determine the nonlinear field relations between the string field theory and effective field theory dilatons and tachyons to quadratic order and at zero-momentum, and examine large dilaton deformations. In the concluding section we discuss additional considerations that suggest the existence of the tachyon vacuum. These come from non-critical string theory, p-adic strings, and sigma model arguments. Finally, the details of the nontrivial computations of quartic couplings are given in the Appendix.
## 2 Computation of the tachyon potential
In this section we present the main computations of this paper. We begin by introducing the string field relevant for the calculation of the tachyon potential, giving a detailed discussion of universality. This string field contains the tachyon, at level zero, the ghost-dilaton, at level two, and massive fields at higher even levels. We then give the quadratic and cubic couplings for the string field restricted to level four and calculate the critical point. Finally, we give the quartic couplings at level zero, two, and four. The critical point survives the inclusion of quartic interactions and becomes more shallow – consistent with the conjecture that the tachyon vacuum has zero action.
The computations use the closed string field action , which takes the form
$$S=\frac{2}{\alpha ^{}}\left(\frac{1}{2}\mathrm{\Psi }|c_0^{}Q|\mathrm{\Psi }+\frac{\kappa }{3!}\{\mathrm{\Psi },\mathrm{\Psi },\mathrm{\Psi }\}+\frac{\kappa ^2}{4!}\{\mathrm{\Psi },\mathrm{\Psi },\mathrm{\Psi },\mathrm{\Psi }\}+\mathrm{}\right).$$
(2.1)
The string field $`\mathrm{\Psi }`$ lives on $``$, the ghost number two state space of the full CFT restricted to the subspace of states that satisfy
$$(L_0\overline{L}_0)|\mathrm{\Psi }=0\text{and}(b_0\overline{b}_0)|\mathrm{\Psi }=0.$$
(2.2)
The BRST operator is $`Q=c_0L_0+\overline{c}_0\overline{L}_0+\mathrm{}`$, where the dots denote terms independent of $`c_0`$ and of $`\overline{c}_0`$. Moreover, $`c_0^\pm =\frac{1}{2}(c_0\pm \overline{c}_0)`$, and we normalize correlators using $`0|c_1\overline{c}_1c_0^{}c_0^+c_1\overline{c}_1|0=1`$. All spacetime coordinates are imagined compactified with the volume of spacetime set equal to one.
### 2.1 Tachyon potential universality and the ghost-dilaton
The universality of the closed string tachyon potential was briefly discussed in , where it was also noted that the ghost number two universal string field that contains the tachyon should include the zero-momentum ghost-dilaton state $`(c_1c_1\overline{c}_1\overline{c}_1)|0`$. In here we review the universality argument and extend it slightly, offering the following observations:
* The ghost-dilaton must be included because closed string field theory is not cubic.
* A world-sheet parity symmetry of closed string field theory can be used to restrict the universal subspace.
* The arguments of do not apply directly to general CFT’s, linear dilaton backgrounds, for example. If the closed string background is defined by a general matter CFT, solutions on the universal subspace may still be solutions, but there is no tachyon potential .
The original idea in universality is to produce a subdivision of all the component fields of the string field theory into two disjoint sets, a set $`\{t_i\}`$ that contains the zero-momentum tachyon and a set $`\{u_a\}`$ such that the string field action $`S(t_i,u_a)`$ contains no term with a single $`u`$-type field. It is then consistent to search for a solution of the equations of motion that assumes $`u_a=0`$ for all $`a`$.
To produce the desired set $`\{t_i\}`$ we assume that the matter CFT is such that $`X^0`$ is the usual negative-metric field with associated conserved momentum $`k_0`$ and the rest of the matter CFT is unitary. The state space $``$ (see (2.2)) is then divided into three disjoint vector subspaces $`_1,_2,`$ and $`_3`$. One has $`_i=_i|𝒢`$, where $`|𝒢`$ denotes a state built with ghost and antighost oscillators only and $`_1,_2,`$ and $`_3`$ are disjoint subspaces of the matter CFT whose union gives the total matter CFT state space:
$`_1:`$ $`\text{the }SL(2,C)\text{ vacuum}|0\text{and descendents},`$
$`_2:`$ $`\text{states with }k_00,`$ (2.3)
$`_3:`$ $`\text{primaries with}k_0=0\text{but different from}|0\text{and descendents}.`$
In the above, primary and descendent refers to the matter Virasoro operators. Note that the primaries in $`_3`$ have positive conformal dimension. The BRST operator preserves the conditions (2.2), and since it is composed of ghost oscillators and matter Virasoro operators, it maps each $`_i`$ into itself. Finally, the spaces $`_i`$ are orthogonal under the BPZ inner product; they only couple to themselves.
The claim is that the set $`\{t_i\}`$ is in fact $`_1`$, the states built upon the zero momentum vacuum. The “tachyon potential” is the string action evaluated for $`_1`$.
We first note that because of momentum conservation fields in $`_2`$ cannot couple linearly to fields in $`_1`$. The fields in $`_3`$ cannot couple linearly to the fields in $`_1`$ either. They cannot do so through the kinetic term because the BRST operator preserves the space and $`_1`$ and $`_3`$ are BPZ orthogonal. We also note that the matter correlator in the $`n`$-string vertex does not couple $`n1`$ vacua $`|0`$ from $`_1`$ to a matter primary from $`_3`$: this is just the one-point function of the primary in $`_3`$, which vanishes because the state has non-zero dimension. The (matter) Virasoro conservation laws on the vertex then imply that the coupling of any $`(n1)`$ states in $`_1`$ to a state in $`_3`$ must vanish. This completes the proof that $`_1`$ is the subspace for tachyon condensation.
The space $`_1`$ can be written as
$$\text{Span}\left\{L_{j_1}^m\mathrm{}L_{j_p}^m\overline{L}_{\overline{j}_1}^m\mathrm{}\overline{L}_{\overline{j}_{\overline{p}}}^mb_{k_1}\mathrm{}b_{k_q}\overline{b}_{\overline{k}_1}\mathrm{}\overline{b}_{\overline{k}_{\overline{q}}}c_{l_1}\mathrm{}c_{l_r}\overline{c}_{\overline{l}_1}\mathrm{}\overline{c}_{\overline{l}_{\overline{r}}}|0\right\},$$
(2.4)
where
$$j_1j_2\mathrm{}j_p,j_i2,\overline{j}_1\overline{j}_2\mathrm{}\overline{j}_{\overline{p}},\overline{j}_i2,$$
(2.5)
as well as
$$k_i,\overline{k}_i2,l_i,\overline{l}_i1,\text{and}r+\overline{r}q\overline{q}=2.$$
(2.6)
Finally, the states above must also be annihilated by $`L_0\overline{L}_0`$ as well as $`b_0\overline{b}_0`$.
There is a reality condition on the string field : its BPZ and hermitian conjugates must differ by a sign. We show now that this condition is satisfied by all the states in (2.4), so the coefficients by which they are multiplied in the universal string field (the zero-momentum spacetime fields) must be real. Suppose a state is built with $`p`$ ghost oscillators and $`p2`$ antighost oscillators. The BPZ and hermitian conjugates differ by the product of two factors: a $`(1)^p`$ from the BPZ conjugation of the ghost oscillators and a $`(1)^{(2p2)(2p1)/2}=(1)^{p1}`$ from the reordering of oscillators in the hermitian conjugate. The product of these two factors is minus one, as we wanted to show.
In open string theory twist symmetry, which arises from world-sheet parity, can be used to further restrict the universal subspace constructed from matter Virasoro and ghost oscillators. In the case of closed string theory the world-sheet parity transformation that exchanges holomorphic and antiholomorphic sectors is the relevant symmetry.<sup>3</sup><sup>3</sup>3We thank A. Sen for discussions that led us to construct the arguments presented below. World-sheet parity is not necessarily a symmetry of arbitrary matter CFT’s, but it is a symmetry in the universal subspace: correlators are complex conjugated when we exchange holomorphic and antiholomorphic Virasoro operators as $`T(z)\overline{T}(\overline{z})`$. More precisely, we introduce a $``$-conjugation, a map of $`_1`$ to $`_1`$ that is an involution. In a basis of Virasoro modes $``$ can be written explicitly as the map of states
$$:AL_{i_1}\mathrm{}L_{i_n}\overline{L}_{j_1}\mathrm{}\overline{L}_{j_n}|0A^{}\overline{L}_{i_1}\mathrm{}\overline{L}_{i_n}L_{j_1}\mathrm{}L_{j_n}|0,$$
(2.7)
where $`A`$ is a constant and $`A^{}`$ denotes its complex conjugate. Given the operator/state correspondence, the above defines completely the star operation $`:𝒪𝒪^{}`$ on vertex operators for vacuum descendents. It results in the following property for the correlator of $`n`$ such operators placed at $`n`$ points on a Riemann surface:
$$𝒪_1\mathrm{}𝒪_n=𝒪_1^{}\mathrm{}𝒪_n^{}^{}.$$
(2.8)
In the ghost sector of the CFT a small complication with signs arises because the basic correlator is odd under the exchange of holomorphic and anti-holomorphic sectors:
$$c(z_1)c(z_2)c(z_3)\overline{c}(\overline{w}_1)\overline{c}(\overline{w}_2)\overline{c}(\overline{w}_3)=\overline{c}(\overline{z}_1)\overline{c}(\overline{z}_2)\overline{c}(\overline{z}_3)c(w_1)c(w_2)c(w_3)^{}.$$
(2.9)
Since two-point functions of the ghost fields are complex conjugated by the exchanges $`c(z)\overline{c}(\overline{z})`$ and $`b(z)\overline{b}(\overline{z})`$, it follows from (2.9) that performing these exchanges on an arbitrary correlator of ghost and antighost fields will give minus the complex conjugate of the original correlator. We will define $``$-conjugation in the ghost sector by:
$$:Ac_{i_1}c_{i_n}b_{j_1}b_{j_m}\overline{c}_{k_1}\overline{c}_{k_r}\overline{b}_{l_1}\overline{b}_{l_s}|0A^{}\overline{c}_{i_1}\overline{c}_{i_n}\overline{b}_{j_1}\overline{b}_{j_m}c_{k_1}c_{k_r}b_{l_1}b_{l_s}|0.$$
(2.10)
For a general state $`\mathrm{\Psi }`$ of the universal subspace we define $`\mathrm{\Psi }^{}`$ to be the state obtained by the simultaneous application of (2.7) and (2.10). It is clear from the above discussion that the correlators satisfy
$$\mathrm{\Psi }_1\mathrm{\Psi }_2\mathrm{}\mathrm{\Psi }_n=\mathrm{\Psi }_1^{}\mathrm{\Psi }_2^{}\mathrm{}\mathrm{\Psi }_n^{}^{},\mathrm{\Psi }_i_1.$$
(2.11)
We now define the action of the world-sheet parity operation $`𝒫`$ on arbitrary states of the universal subspace:
$$𝒫\mathrm{\Psi }\mathrm{\Psi }^{},\mathrm{\Psi }_1.$$
(2.12)
We claim that the string field theory action, restricted to $`_1`$, is $`𝒫`$ invariant:
$$S(\mathrm{\Psi })=S(𝒫\mathrm{\Psi }),\text{for}\mathrm{\Psi }_1.$$
(2.13)
First consider the invariance of the cubic term. Using (2.12) and (2.11) we have
$$𝒫\mathrm{\Psi },𝒫\mathrm{\Psi },𝒫\mathrm{\Psi }=\mathrm{\Psi }^{},\mathrm{\Psi }^{},\mathrm{\Psi }^{}=\mathrm{\Psi },\mathrm{\Psi },\mathrm{\Psi }^{}=\mathrm{\Psi },\mathrm{\Psi },\mathrm{\Psi },$$
(2.14)
where in the last step we used the reality of the string field action. The kinetic term of the action is also invariant. First note that $`(c_0^{}Q\mathrm{\Psi })^{}=c_0^{}Q\mathrm{\Psi }^{}.`$ It then follows that
$$𝒫\mathrm{\Psi },c_0^{}Q𝒫\mathrm{\Psi }=\mathrm{\Psi }^{},c_0^{}Q\mathrm{\Psi }^{}=\mathrm{\Psi }^{},(c_0^{}Q\mathrm{\Psi })^{}=\mathrm{\Psi },c_0^{}Q\mathrm{\Psi }^{}=\mathrm{\Psi },c_0^{}Q\mathrm{\Psi }.$$
(2.15)
For higher point interactions, the invariance follows because the antighost insertions have the appropriate structure. Each time we add a new string field we must add two antighost insertions. For the case of quartic interactions they take the form of two factors $`^{}`$ (see eqn. (A.3)). Since $`(^{})^{}=^{}`$, the extra minus sign cancels against the minus sign from the extra string field. This can be seen to generalize to higher order interactions using the forms of the off-shell amplitudes discussed in section 6 of . This completes our proof of (2.13).
Since $`𝒫^2=1`$ the space $`_1`$ can be divided into two disjoint subspaces: the space $`_1^+`$ of states with $`𝒫=1`$ and the space $`_1^{}`$ of states with $`𝒫=1`$:
$`𝒫(\mathrm{\Psi }_+)`$ $`=`$ $`+\mathrm{\Psi }_+,\mathrm{\Psi }_+_1^+,`$
$`𝒫(\mathrm{\Psi }_{})`$ $`=`$ $`\mathrm{\Psi }_{},\mathrm{\Psi }_{}_1^{}.`$ (2.16)
It follows from the invariance of the action that no term in the action can contain just one state in $`_1^{}`$. We can therefore restrict ourselves to the subspace $`_1^+`$ with positive parity.
The string field is further restricted by using a gauge fixing condition. The computation of the potential is done in the Siegel gauge, which requires states to be annihilated by $`b_0+\overline{b}_0`$. To restrict ourselves to the Siegel gauge we take the states in (2.4) that have neither a $`c_0`$ nor a $`\overline{c}_0`$.
The Siegel gauge fixes the gauge symmetry completely for the massive levels, but does not quite do the job at the massless level. There are two states with $`L_0=\overline{L}_0=0`$ in $`_1`$ that are in the Siegel gauge:
$$(c_1c_1\overline{c}_1\overline{c}_1)|0\text{and}(c_1c_1+\overline{c}_1\overline{c}_1)|0.$$
(2.17)
The first state is the ghost dilaton and it is proportional to $`Q(c_0\overline{c}_0)|0`$. Since $`(c_0\overline{c}_0)|0`$ is not annihilated by $`b_0\overline{b}_0`$ the gauge parameter is illegal and the ghost dilaton is not trivial. The second state is proportional to $`Q(c_0+\overline{c}_0)|0`$, so it is thus trivial at the linearized level. Although trivial at the linearized level, one may wonder if the triviality holds for large fields. Happily, we need not worry: the state is $`𝒫`$ odd, so it need not be included in the calculation. The ghost-dilaton, because of the relative minus sign between the two terms, is $`𝒫`$ even and it is included.
Had the closed string field theory been cubic we could have discarded the ghost-dilaton state and all other states with asymmetric left and right ghost numbers. We could restrict $`_1^+`$ to fields of ghost number $`(G,\overline{G})=(1,1)`$. Indeed, the cubic vertex cannot couple two $`(1,1)`$ fields to anything except another $`(1,1)`$ field. Moreover, in the Siegel gauge $`c_0^{}Q`$ acts as an operator of ghost number $`(1,1)`$, so again, no field with asymmetric ghost numbers can couple linearly. The quartic and higher order interactions in CSFT have antighost insertions that do not have equal left and right ghost numbers. It follows that these higher order vertices can couple the ghost-dilaton to $`(1,1)`$ fields. Indeed, the coupling of a dilaton to three tachyons does not vanish. We cannot remove from $`_1^+`$ the dilaton, nor other states with asymmetric left and right ghost numbers.
The construction of the universal string field and action presented here does not work fully if the matter CFT contains a linear dilaton background. Momentum conservation along the corresponding coordinate is anomalous and one cannot build an action with states of zero momentum only: the action restricted to $`_1`$ is identically zero. There would be no universal “potential” in $`_1`$. It appears rather likely, however, that any solution in the universal subspace would still be a solution in a linear dilaton background. In fact, any solution in the universal subspace may be a solution for string field theory formulated with a general matter CFT .
We conclude this section by writing out the string field for the first few levels. The level $`\mathrm{}`$ of a state is defined by $`\mathrm{}=L_0+\overline{L}_0+2.`$ The level zero part of the string field is
$$|\mathrm{\Psi }_0=tc_1\overline{c}_1|0.$$
(2.18)
Here $`t`$ is the zero-momentum tachyon. The level two part of the string field is
$$|\mathrm{\Psi }_2=d(c_1c_1\overline{c}_1\overline{c}_1)|0.$$
(2.19)
Here $`d`$ is the zero momentum ghost-dilaton. It multiplies the only state of $`𝒫=+1`$ at this level. At level four there are four component fields:
$`|\mathrm{\Psi }_4`$ $`=`$ $`(f_1c_1\overline{c}_1+f_2L_2c_1\overline{L}_2\overline{c}_1+f_3(L_2c_1\overline{c}_1+c_1\overline{L}_2\overline{c}_1)`$ (2.20)
$`+g_1(b_2c_1\overline{c}_2\overline{c}_1c_2c_1\overline{b}_2\overline{c}_1))|0.`$
Note that the states coupling to the component fields all have $`𝒫=+1`$ and that $`g_1`$ couples to a state with asymmetric left and right ghost numbers. In this paper we will not use higher level terms in the string field.
With $`\alpha ^{}=2`$ the closed string field potential $`V`$ associated with the action in (2.1) is
$$\kappa ^2V=\frac{1}{2}\mathrm{\Psi }|c_0^{}Q|\mathrm{\Psi }+\frac{1}{3!}\{\mathrm{\Psi },\mathrm{\Psi },\mathrm{\Psi }\}+\frac{1}{4!}\{\mathrm{\Psi },\mathrm{\Psi },\mathrm{\Psi },\mathrm{\Psi }\}+\mathrm{}.$$
(2.21)
Here $`|\mathrm{\Psi }=|\mathrm{\Psi }_0+|\mathrm{\Psi }_2+|\mathrm{\Psi }_4+\mathrm{}`$. Our computations will not include quintic and higher order interactions in the string action.
### 2.2 The quadratic and cubic terms in the potential
Let us now consider the potential including only the kinetic and cubic terms in (2.21). To level zero:
$$\kappa ^2V_0^{(2)}=t^2,\kappa ^2V_0^{(3)}=\frac{6561}{4096}t^3.$$
(2.22)
All potentials introduced in this subsection have a superscript that gives the order of the interaction (two for quadratic, three for cubic, and so on), and a subscript that gives the level (defined by the sum of levels of fields in the interaction). The next terms arise at level four, where we have couplings of the tachyon to the square of the dilaton and couplings of the level four fields to the tachyon squared:
$$\kappa ^2V_4^{(3)}=\frac{27}{32}d^2t+\left(\frac{3267}{4096}f_1+\frac{114075}{4096}f_2\frac{19305}{2048}f_3\right)t^2.$$
(2.23)
At level six we can couple a level four field, a dilaton, and a tachyon. Only level four fields with $`G\overline{G}`$ can have such coupling, so we find:
$$\kappa ^2V_6^{(3)}=\frac{25}{8}g_1td.$$
(2.24)
At level eight there are two kinds of terms. First, we have the kinetic terms for the level four fields:
$$\kappa ^2V_8^{(2)}=f_{1}^{}{}_{}{}^{2}+169f_{2}^{}{}_{}{}^{2}26f_{3}^{}{}_{}{}^{2}2g_{1}^{}{}_{}{}^{2}.$$
(2.25)
Second, we have the cubic interactions:
$$\begin{array}{cc}\hfill \kappa ^2V_8^{(3)}& =\frac{1}{96}f_1d^2\frac{4225}{864}f_2d^2+\frac{65}{144}f_3d^2\hfill \\ & +\frac{361}{12288}f_{1}^{}{}_{}{}^{2}t+\frac{511225}{55296}f_1f_2t+\frac{57047809}{110592}f_{2}^{}{}_{}{}^{2}t+\frac{470873}{27648}f_{3}^{}{}_{}{}^{2}t\frac{49}{24}g_{1}^{}{}_{}{}^{2}t\hfill \\ & \frac{13585}{9216}f_1f_3t\frac{5400395}{27648}f_2f_3t.\hfill \end{array}$$
(2.26)
As we can see, these are of two types: couplings of a level four field to two dilatons (first line) and couplings of two level four fields to a tachyon (second and third lines).
The terms at level 10 couple two level four fields and a dilaton. Because of ghost number conservation, one of the level four fields must have $`G\overline{G}`$:
$$\kappa ^2V_{10}^{(3)}=\frac{25}{5832}\left(361f_1+4225f_22470f_3\right)dg_1.$$
(2.27)
Finally, at level 12 we have the cubic couplings of three level-four fields:
$$\begin{array}{cc}\hfill \kappa ^2V_{12}^{(3)}& =\frac{1}{4096}f_1^3+\frac{1525225}{8957952}f_1^2f_2\frac{1235}{55296}f_1^2f_3+\frac{6902784889}{80621568}f_1f_2^2\hfill \\ & \frac{102607505}{6718464}f_1f_2f_3+\frac{1884233}{2239488}f_1f_3^2\hfill \\ & +\frac{74181603769}{26873856}f_2^3\frac{22628735129}{13436928}f_2^2f_3+\frac{4965049817}{20155392}f_2f_3^2\hfill \\ & \frac{31167227}{3359232}f_3^3\frac{961}{157464}f_1g_1^2\frac{207025}{17496}f_2g_1^2+\frac{14105}{26244}f_3g_1^2.\hfill \end{array}$$
(2.28)
### 2.3 Tachyon vacuum with cubic vertices only
With cubic vertices only the dilaton expectation value is zero. In fact, only fields with $`G=\overline{G}=1`$ can acquire nonvanishing expectation values. To examine the tachyon vacuum we define a series of potentials:
$$\begin{array}{cc}\hfill 𝕍_0^{(3)}& V_0^{(2)}+V_0^{(3)},\hfill \\ \hfill 𝕍_8^{(3)}& 𝕍_0^{(3)}+V_4^{(3)}+V_6^{(3)}+V_8^{(2)}+V_8^{(3)},\hfill \\ \hfill 𝕍_{12}^{(3)}& 𝕍_8^{(3)}+V_{10}^{(3)}+V_{12}^{(3)}.\hfill \end{array}$$
(2.29)
A few observations are in order. In all of the above potentials we can set $`d=g_1=0`$. As a consequence, $`V_6^{(3)}`$ and $`V_{10}^{(3)}`$ do not contribute. Since the level-two dilaton plays no role, once we go beyond the tachyon we must include level four fields. The kinetic terms for these fields are of level eight, so $`𝕍_8^{(3)}`$ is the simplest potential beyond level zero. With level-four fields the next potential is $`𝕍_{12}^{(3)}`$.
The critical points obtained with the potentials $`𝕍_0^{(3)}`$, $`𝕍_8^{(3)}`$, and $`𝕍_{12}^{(3)}`$ are given in Table 1. We call the value of the potential $`\kappa ^2𝕍`$ at the critical point the action density. The values of the action density follow the pattern of open string theory. The original cubic critical point becomes deeper. It does so by about 10%, a value significantly smaller than the corresponding one in open string field theory.
### 2.4 Tachyon vacuum with cubic and quartic vertices
We can now examine the quartic terms in the potential. The associated potentials are denoted with a superscript $`(4)`$ for quartic and a subscript that gives the sum of levels of the fields that enter the term. The quartic self-coupling of tachyons has been calculated in :
$$\kappa ^2V_0^{(4)}=3.0172t^4.$$
(2.30)
With total level two we have a coupling of three tachyons and one dilaton. This is calculated in Appendix A.2 and the result is
$$\kappa ^2V_2^{(4)}=\mathrm{\hspace{0.17em}3.8721}t^3d.$$
(2.31)
With total level four there is the coupling of two tachyons to two dilatons (Appendix A.2) and the coupling of three tachyons to any of the level-four fields (Appendix A.3):
$$\begin{array}{c}\hfill \kappa ^2V_4^{(4)}=1.3682t^2d^2+t^3\left(0.4377f_156.262f_2+13.024f_3+\mathrm{\hspace{0.17em}0.2725}g_1\right).\end{array}$$
(2.32)
With total level six there are three types of interactions: a tachyon coupled to three dilatons, two tachyons coupled to a dilaton and a level-four field, and three tachyons coupled to a level-six field. We have only computed the first one (Appendix A.2):
$$\kappa ^2V_6^{(4)}=\mathrm{\hspace{0.17em}0.9528}td^3+\mathrm{}.$$
(2.33)
The terms that have not been computed are indicated by the dots. Finally, the quartic self-coupling of dilatons was computed in , where it played a central role in the demonstration that the effective dilaton potential has no quartic term:
$$\kappa ^2V_8^{(4)}=\mathrm{\hspace{0.17em}0.1056}d^4+\mathrm{}.$$
(2.34)
We use the dots to indicate the additional level eight interactions that should be computed.
Let us now consider the potentials that can be assembled using the above contributions. We use the following strategy: we include cubic vertices to the highest possible level and then begin to introduce the quartic couplings level by level. The most accurate potential with quadratic and cubic terms that we have is $`𝕍_{12}^{(3)}`$ and the tachyon vacuum it contains appears in the last line of Table 1. The lowest order quartic potential that we use is therefore:
$$𝕍_0^{(4)}𝕍_{12}^{(3)}+V_0^{(4)}.$$
(2.35)
This potential has a familiar difficulty: the quartic self-coupling of the tachyon is so strong that the critical point in the potential disappears. As we have argued, once additional terms are included the critical point in the potential reappears. The higher level potentials are defined by including progressively higher level quartic interactions:
$$\begin{array}{cc}\hfill 𝕍_2^{(4)}& 𝕍_0^{(4)}+V_2^{(4)},\hfill \\ \hfill 𝕍_4^{(4)}& 𝕍_2^{(4)}+V_4^{(4)}.\hfill \end{array}$$
(2.36)
Since our computations of $`V_6^{(4)}`$ and $`V_8^{(4)}`$ are incomplete, the results that follow from $`𝕍_6^{(4)}𝕍_4^{(4)}+V_6^{(4)}`$ and $`𝕍_8^{(4)}𝕍_6^{(4)}+V_8^{(4)}`$ cannot be trusted.
We are now in a position to calculate the critical points of the potentials $`𝕍^{(4)}`$. In our numerical work we input the cubic coefficients as fractions and the quartic coefficients as the exact decimals given above (so the $`t^4`$ coefficient is treated as exactly equal to $`3.0172`$.) Our results are given in Table 2. For ease of comparison, we have included the cubic results for $`𝕍_{12}^{(3)}`$ as the first line. Furthermore, we include a line for $`𝕍_0^{(4)}`$ even though there is no critical point. The next potential is $`𝕍_2^{(4)}`$ which contains only the additional coupling $`t^3d`$. The significant result is that the critical point reappears and can be considered to be a (moderate) deformation of the critical point obtained with $`𝕍_{12}^{(3)}`$. Indeed, while there is a new expectation value for the dilaton (and for $`g_1`$), the expectation value of the tachyon does not change dramatically, nor do the expectation values for $`f_1`$, $`f_2`$, and $`f_3`$. The critical point becomes somewhat shallower, despite the destabilizing effects of the tachyon quartic self-couplings.
At the next level, where $`t^2d^2`$ and $`t^3M_4`$ ($`M_4`$ denotes a level-four field) terms appear, the critical point experiences some significant change. First of all, it becomes about 40% more shallow; the change is large and probably significant, given the expectation that the action density should eventually reach zero. The tachyon expectation changes considerably but the dilaton expectation value changes little. Due to the $`t^3M_4`$ terms the expectation values of some of the level four fields change dramatically.
Glancing at Table 2, one notices that the tachyon expectation value is becoming smaller so one might worry that the critical point is approaching the perturbative vacuum. This is, of course, a possibility. If realized, it would imply that the critical point we have encountered is an artifact of level expansion. We think this is unlikely. Since the dilaton seems to be relatively stable, a trivial critical point would have to be a dilaton deformation of the perturbative vacuum, but such deformations have negative tachyon expectation values (see Figure 2).
At this moment we do not have full results for higher levels. The computation of $`𝕍_6^{(4)}`$ would require the evaluation of couplings of the form $`t^2dM_4`$ and, in principle, couplings $`t^3M_6`$ of level-six fields, which we have not even introduced in this paper. The only additional couplings we know at present are $`td^3`$, which enters in $`𝕍_6^{(4)}`$ and $`d^4`$, which enters in $`𝕍_8^{(4)}`$ (see eqns. (2.33) and (2.34)). Despite lacking terms, we calculated the resulting vacua to test that no wild effects take place. The incomplete $`𝕍_6^{(4)}`$ leads to $`t=0.35426,d=0.40763`$ and an action density of $`0.05553`$. The incomplete $`𝕍_8^{(4)}`$ leads to $`t=0.36853,d=0.40222`$ and an action density of $`0.05836`$. In these results the action density has become more negative. Given the conjectured value of the action, it would be encouraging if the full results at those levels show an action density whose magnitude does not become larger.
One may also wonder what happens if terms of order higher than quartic are included in the potential. Since the tachyon terms in the CSFT potential alternate signs , the quintic term is positive and will help reduce the value of the action at the critical point. The coefficient of this coupling will be eventually needed as computations become more accurate. The sixtic term will have a destabilizing effect. Having survived the destabilizing effects of the quartic term, we can hope that those of the sixtic term will prove harmless. If, in general, even power terms do not have catastrophic effects, it may be better to work always with truncations of odd power.
## 3 The sigma model and the string field theory pictures
In this section we study the relations between the string field metric $`h_{\mu \nu }`$ and the ghost-dilaton $`d`$ and the corresponding sigma model fields, the string metric $`\stackrel{~}{h}_{\mu \nu }`$ and dilaton $`\mathrm{\Phi }`$. These relations are needed to interpret the tachyon vacuum solution and to discuss the possible relation to the rolling solutions.
We begin by finding the precise linearized relations between the string field dilaton and the sigma model dilaton. The linearized relations confirm that the CSFT metric $`h_{\mu \nu }`$, which does not acquire an expectation value in the tachyon vacuum, coincides with the string metric of the sigma model, which does not change in the rolling solutions. Moreover, the relation (3.14), together with $`h_{\mu \nu }=0`$, implies that our $`d>0`$ in the tachyon vacuum corresponds to $`\mathrm{\Phi }>0`$, thus larger string coupling. This is also consistent with what we obtained in the rolling solutions.
Our discussion of the linearized relations also allows us to examine the various vertex operators associated with the various dilaton fields used in the literature (section 3.2.). In section 3.3 we examine the nonlinear relations between the CSFT tachyon and dilaton and the effective field theory ones. We work at zero momentum and up to quadratic order. Finally, in section 3.4, we present evidence that CSFT can describe arbitrarily large dilaton deformations.
### 3.1 Relating sigma model fields and string fields
Consider first the effective action (1.5), suggested by the conditions of conformal invariance of a sigma model with gravity, dilaton and tachyon background fields. If we set the tachyon to zero, this action reduces to the effective action for massless fields, in the conventions of . In this action $`g_{\mu \nu }`$ is the string metric, $`\mathrm{\Phi }`$ is the diffeomorphism invariant dilaton, and $`T`$, with potential $`V(T)=\frac{2}{\alpha ^{}}T^2+\mathrm{}`$, is the tachyon. In order to compare with the string field action we expand the effective action in powers of small fluctuations using
$$g_{\mu \nu }=\eta _{\mu \nu }+\stackrel{~}{h}_{\mu \nu },$$
(3.1)
where we use a tilde in the fluctuation to distinguish it from the metric fluctuation in the string field. The result is
$$\begin{array}{cc}\hfill S_\sigma & =\frac{1}{2\kappa ^2}d^Dx(\frac{1}{4}\stackrel{~}{h}_{\mu \nu }^2\stackrel{~}{h}^{\mu \nu }\frac{1}{4}\stackrel{~}{h}^2\stackrel{~}{h}+\frac{1}{2}(^\nu \stackrel{~}{h}_{\mu \nu })^2+\frac{1}{2}\stackrel{~}{h}_\mu _\nu \stackrel{~}{h}^{\mu \nu }\hfill \\ & +2\stackrel{~}{h}^2\mathrm{\Phi }2\mathrm{\Phi }_\mu _\nu \stackrel{~}{h}^{\mu \nu }4\mathrm{\Phi }^2\mathrm{\Phi }\hfill \\ & (T)^2+\frac{4}{\alpha ^{}}T^2+\stackrel{~}{h}^{\mu \nu }_\mu T_\nu T+(\frac{\stackrel{~}{h}}{2}2\mathrm{\Phi })(T)^2+\mathrm{}),\hfill \end{array}$$
(3.2)
where we have kept cubic terms coupling the dilaton and metric to the tachyon. Such terms are needed to fix signs in the relations between the fields in the sigma model and the string fields.
Let us now consider the string field action. The string field needed to describe the tachyon, the metric fluctuations, and the dilaton is
$$\begin{array}{cc}\hfill |\mathrm{\Psi }& =\frac{d^Dk}{(2\pi )^D}(t(k)c_1\overline{c}_1\frac{1}{2}h_{\mu \nu }(k)\alpha _1^\mu \overline{\alpha }_1^\nu c_1\overline{c}_1+d(k)(c_1c_1\overline{c}_1\overline{c}_1)\hfill \\ & +i\sqrt{\frac{\alpha ^{}}{2}}B_\mu (k)c_0^+(c_1\alpha _1^\mu \overline{c}_1\overline{\alpha }_1^\mu ))|k.\hfill \end{array}$$
(3.3)
Here $`t(k)`$ is the tachyon, $`h_{\mu \nu }(k)=h_{\nu \mu }(k)`$ is a metric fluctuation, $`d(k)`$ is the ghost-dilaton, and $`B_\mu (k)`$ is an auxiliary field. The sign and coefficient of $`h_{\mu \nu }`$ have been chosen for future convenience. The linearized gauge transformations of the component fields can be obtained from $`\delta |\mathrm{\Psi }=Q_B|\mathrm{\Lambda }`$ with
$$|\mathrm{\Lambda }=\frac{i}{\sqrt{2\alpha ^{}}}ϵ_\mu (c_1\alpha _1^\mu \overline{c}_1\overline{\alpha }_1^\mu )|p.$$
(3.4)
The resulting coordinate-space gauge transformations are:
$$\delta h_{\mu \nu }=_\nu ϵ_\mu +_\mu ϵ_\nu ,\delta d=\frac{1}{2}ϵ,\delta B_\mu =\frac{1}{2}^2ϵ_\mu ,\delta t=0.$$
(3.5)
We now calculate the quadratic part of the closed string field action, finding
$$\begin{array}{cc}\hfill S^{(2)}& =\frac{1}{\kappa ^2\alpha ^{}}\mathrm{\Psi }|c_0^{}Q_B|\mathrm{\Psi },\hfill \\ & =\frac{1}{2\kappa ^2}d^Dx\left(\frac{1}{4}h_{\mu \nu }^2h^{\mu \nu }2d^2d2B_\mu (_\nu h^{\mu \nu }+2^\mu d)2B^2(t)^2+\frac{4}{\alpha ^{}}t^2\right),\hfill \\ & =\frac{1}{2\kappa ^2}d^Dx\left(\frac{1}{4}h_{\mu \nu }^2h^{\mu \nu }+\frac{1}{2}(^\nu h_{\mu \nu })^24d^2d2d_\mu _\nu h^{\mu \nu }(t)^2+\frac{4}{\alpha ^{}}t^2\right).\hfill \end{array}$$
(3.6)
In the last step we eliminated the auxiliary field $`B_\mu `$ using its algebraic equation of motion.
The gauge transformations (3.5) imply that the linear combination $`d+\frac{h}{4}`$ is gauge invariant. It follows that the sigma model dilaton must take the form
$$\lambda \mathrm{\Phi }=d+\frac{h}{4},$$
(3.7)
where $`\lambda `$ is a number to be determined. Using (3.7) to eliminate the ghost-dilaton $`d`$ from the action (3.6) we find
$$\begin{array}{cc}\hfill S^{(2)}& =\frac{1}{2\kappa ^2}d^Dx(\frac{1}{4}h_{\mu \nu }^2h^{\mu \nu }\frac{1}{4}h^2h+\frac{1}{2}(^\nu h_{\mu \nu })^2+\frac{1}{2}h_\mu _\nu h^{\mu \nu }\hfill \\ & +\mathrm{\hspace{0.17em}2}\lambda h^2\mathrm{\Phi }2\lambda \mathrm{\Phi }_\mu _\nu h^{\mu \nu }4\lambda ^2\mathrm{\Phi }^2\mathrm{\Phi }(t)^2+\frac{4}{\alpha ^{}}t^2).\hfill \end{array}$$
(3.8)
We also use the string field theory to calculate the on-shell coupling of $`h_{\mu \nu }`$ to two tachyons. This coupling arises from the term
$$S^{(3)}=\frac{1}{\alpha ^{}\kappa ^2}𝒯,,𝒯,$$
(3.9)
where $`𝒯`$ and $``$ denote the parts of the string field (3.3) that contain $`t(k)`$ and $`h_{\mu \nu }(k)`$, respectively. We thus have
$$S^{(3)}=\frac{1}{2\alpha ^{}\kappa ^2}\left(\underset{i=1}{\overset{3}{}}\frac{d^Dk_i}{(2\pi )^D}\right)c_1\overline{c}_1e^{ik_1X},c_1\overline{c}_1\alpha _1^\mu \overline{\alpha }_1^\nu e^{ik_2X},c_1\overline{c}_1e^{ik_3X}t(k_1)t(k_3)h_{\mu \nu }(k_2).$$
(3.10)
The on-shell evaluation is readily carried out using $`k^\mu h_{\mu \nu }(k)=0`$. We obtain
$$S^{(3)}=\frac{1}{2\kappa ^2}\frac{d^Dk_1}{(2\pi )^D}\frac{d^Dk_3}{(2\pi )^D}k_1^\mu k_3^\nu t(k_1)t(k_3)h_{\mu \nu }(k_1k_3)=\frac{1}{2\kappa ^2}d^Dxh^{\mu \nu }_\mu t_\nu t.$$
(3.11)
Combining this result with (3.8) we obtain the closed string field theory action
$$\begin{array}{cc}\hfill S_{csft}& =\frac{1}{2\kappa ^2}d^Dx(\frac{1}{4}h_{\mu \nu }^2h^{\mu \nu }\frac{1}{4}h^2h+\frac{1}{2}(^\nu h_{\mu \nu })^2+\frac{1}{2}h_\mu _\nu h^{\mu \nu }\hfill \\ & +\mathrm{\hspace{0.17em}2}\lambda h^2\mathrm{\Phi }2\lambda \mathrm{\Phi }_\mu _\nu h^{\mu \nu }4\lambda ^2\mathrm{\Phi }^2\mathrm{\Phi }\hfill \\ & (t)^2+\frac{4}{\alpha ^{}}t^2+h^{\mu \nu }_\mu t_\nu t+\mathrm{}).\hfill \end{array}$$
(3.12)
We are finally in a position to identify the sigma model action (3.2) and the string field action (3.12). Comparing the quadratic terms in $`\stackrel{~}{h}_{\mu \nu }`$ and those in $`h_{\mu \nu }`$ we see that $`\stackrel{~}{h}_{\mu \nu }=\pm h_{\mu \nu }`$. We also note that $`T=\pm t`$. The coupling $`\stackrel{~}{h}^{\mu \nu }_\mu T_\nu T`$ in (3.2) coincides with the corresponding coupling in (3.12) if and only if
$$\stackrel{~}{h}_{\mu \nu }=h_{\mu \nu }.$$
(3.13)
This simple equality justifies the multiplicative factor of $`(1/2)`$ introduced for $`h_{\mu \nu }`$ in the string field (3.3). The string field $`h_{\mu \nu }`$ so normalized is the fluctuation of the string metric. Comparing the couplings of metric and dilaton in both actions we also conclude that $`\lambda =+1`$ and, therefore, equation (3.7) gives
$$\mathrm{\Phi }=d+\frac{h}{4}.$$
(3.14)
This expresses the sigma model dilaton $`\mathrm{\Phi }`$ in terms of the string field metric trace and the ghost dilaton $`d`$. It is important to note that when we give a positive expectation value to $`d`$ (and no expectation value to $`h`$) we are increasing the value of $`\mathrm{\Phi }`$ and therefore increasing the value of the string coupling.
### 3.2 The many faces of the dilaton
Equipped with the precise relations between string fields and sigma-model fields we digress on the various dilaton fields used in the literature. Of particular interest are the corresponding vertex operators, which are determined by the CFT states that multiply the component fields in the closed string field.
We introduce the states
$$|𝒪^{\mu \nu }(p)=\frac{1}{4}(\alpha _1^\mu \overline{\alpha }_1^\nu +\alpha _1^\nu \overline{\alpha }_1^\mu )|p,|𝒪^d(p)=(c_1c_1\overline{c}_1\overline{c}_1)|p.$$
(3.15)
The corresponding vertex operators are
$$𝒪^{\mu \nu }(p)=\frac{1}{2\alpha ^{}}(X^\mu \overline{}X^\nu +X^\nu \overline{}X^\mu )e^{ipX},𝒪^d(p)=\frac{1}{2}(c^2c\overline{c}\overline{}^2c)e^{ipX}.$$
(3.16)
Working for fixed momentum, the string field (3.3) restricted to metric and dilaton fluctuations is
$$|\mathrm{\Psi }=h_{\mu \nu }|𝒪^{\mu \nu }+d|𝒪^d.$$
(3.17)
This equation states that $`𝒪^d`$ is the vertex operator associated with the ghost-dilaton field $`d`$. An excitation by this vertex operator does not change the metric $`h_{\mu \nu }`$. Our transformation to a gauge invariant dilaton gives
$$\mathrm{\Phi }=d+\frac{1}{4}h,\stackrel{~}{h}_{\mu \nu }=h_{\mu \nu }.$$
(3.18)
Here $`\stackrel{~}{h}_{\mu \nu }`$ is the fluctuation of the string metric. Inverting these relations
$$d=\mathrm{\Phi }\frac{1}{4}\stackrel{~}{h},h_{\mu \nu }=\stackrel{~}{h}_{\mu \nu }.$$
(3.19)
Subtituting into the string field (3.17) we obtain
$$|\mathrm{\Psi }=\stackrel{~}{h}_{\mu \nu }\left(|𝒪^{\mu \nu }\frac{1}{4}\eta ^{\mu \nu }|𝒪^d\right)+\mathrm{\Phi }|𝒪^d.$$
(3.20)
It is interesting to note that $`𝒪^d`$ is the vertex operator associated with a variation of the gauge-invariant dilaton $`\mathrm{\Phi }`$ and no variation of the string metric. On the other hand, $`𝒪^{\mu \nu }\frac{1}{4}\eta ^{\mu \nu }𝒪^d`$ varies the string metric and does not vary the gauge-invariant dilaton (although it varies the ghost-dilaton).
Finally, we consider the formulation that uses the Einstein metric $`g_{\mu \nu }^E`$ and the dilaton $`\mathrm{\Phi }`$. The field redefinition is
$$g_{\mu \nu }^E=\mathrm{exp}(2\omega )g_{\mu \nu },\text{with}\omega =\frac{2}{D2}\mathrm{\Phi }.$$
(3.21)
Expanding in fluctuation fields we obtain
$$h_{\mu \nu }^E=\stackrel{~}{h}_{\mu \nu }\frac{4}{D2}\eta _{\mu \nu }\mathrm{\Phi }.$$
(3.22)
Solving for $`d`$ and $`h_\mu `$ in terms of $`\mathrm{\Phi }`$ and $`h_{\mu \nu }^E`$ we get
$$d=\frac{2}{D2}\mathrm{\Phi }\frac{1}{4}h^E,h_{\mu \nu }=h_{\mu \nu }^E+\frac{4}{D2}\eta _{\mu \nu }\mathrm{\Phi }.$$
(3.23)
Substituting into the string field (3.17) we obtain
$$|\mathrm{\Psi }=h_{\mu \nu }^E\left(|𝒪^{\mu \nu }\frac{1}{4}\eta ^{\mu \nu }|𝒪^d\right)+\frac{2}{D2}\mathrm{\Phi }\left(\mathrm{\hspace{0.17em}2}\eta _{\mu \nu }|𝒪^{\mu \nu }|𝒪^d\right).$$
(3.24)
Interestingly, the vertex operator that varies the Einstein metric (without variation of the dilaton) is the same as that for the string metric (see (3.20)). It is the dilaton operator that changes this time. The vertex operator
$$𝒟=2\eta _{\mu \nu }𝒪^{\mu \nu }𝒪^d=\left(\frac{2}{\alpha ^{}}X\overline{}X\frac{1}{2}(c^2c\overline{c}\overline{}^2c)\right)e^{ipX},$$
(3.25)
varies the dilaton without varying the Einstein metric. This is the dilaton vertex operator used almost exclusively in the early literature – it is naturally associated with the Einstein metric. The corresponding state $`|𝒟(p)`$ has a particularly nice property: it is annihilated by the BRST operator when $`p^2=0`$. Indeed,
$$Q_B|𝒟(p)=\frac{\alpha ^{}}{2}p^2c_0^+|𝒟(p).$$
(3.26)
The dilaton $`𝒟`$ is in fact the unique linear combination of the matter and ghost dilatons that has this property. For other combinations, terms linear in the momentum $`p`$ (such as $`(p\alpha _1)c_1\overline{c}_1\overline{c}_1|p`$), survive.
### 3.3 Relating the sigma model and string field dilaton and tachyon
The closed string theory potential $`V`$, as read from the effective action (1.5) is
$$\kappa ^2V=e^{2\mathrm{\Phi }}\left(V(T)+\mathrm{}\right),\text{with}V(T)=T^2+\mathrm{}.$$
(3.27)
Here $`\mathrm{\Phi }`$ and $`T`$ are the zero momentum dilaton and tachyon fields in the effective field theory. The purpose of this section is to discuss the relation between $`\mathrm{\Phi }`$ and $`T`$ and the corresponding string fields $`d`$ and $`t`$, both sets at zero-momentum. To do this we must consider the effective potential for $`d`$ and $`t`$ calculated in string field theory. We only have the potential itself. Collecting our previous results, we write
$`\kappa ^2V`$ $`=`$ $`t^2+1.6018t^33.0172t^4`$ (3.28)
$`+\mathrm{\hspace{0.17em}3.8721}t^3d+(0.8438t+1.3682t^2)d^20.9528td^30.1056d^4.`$
The contributions from massive fields affect quartic and higher order terms. In our setup, the relevant terms arise when we eliminate the level-four massive fields using their kinetic terms in (2.25) and their linear couplings to $`t^2`$ in (2.23), to $`td`$ in (2.24), and to $`d^2`$ in (2.26). We find
$$\mathrm{\Delta }V=\frac{6241}{186624}d^4+\frac{25329}{16384}d^2t^2\frac{1896129}{4194304}t^40.0334d^4+1.5460d^2t^20.4521t^4.$$
(3.29)
It follows that the effective potential for the tachyon and the dilaton, calculated up to terms quartic in the fields and including massive fields of level four only, is given by:
$`\kappa ^2V_{eff}`$ $`=`$ $`t^2+1.6018t^33.4693t^4`$ (3.30)
$`+\mathrm{\hspace{0.17em}3.8721}t^3d+(0.8438t+2.9142t^2)d^20.9528td^30.1390d^4+\mathrm{}.`$
The dots represent quintic and higher terms, which receive contributions both from elementary interactions and some integration of massive fields. We write, more generically
$`\kappa ^2V_{eff}`$ $`=`$ $`t^2+a_{3,0}t^3+a_{4,0}t^4`$ (3.31)
$`+a_{3,1}t^3d+(a_{1,2}t+a_{2,2}t^2)d^2+a_{1,3}td^3+a_{0,4}d^4+\mathrm{}.`$
The values of the coefficients $`a_{i,j}`$ can be read comparing this equation with (3.30).
There are two facts about $`V_{eff}`$ that make it clear it is not in the form of a ghost-dilaton exponential times a tachyon potential. First, it does not have a term of the form $`t^2d`$ that would arise from the tachyon mass term and the expansion of the exponential. Second, it contains a term linear in the tachyon; those terms should be absent since the tachyon potential does not have a linear term. Nontrivial field redefinitions are necessary to relate string fields and sigma model fields.
To linearized order the fields are the same, so we write relations of the form :
$`t`$ $`=`$ $`T+\alpha _1T\mathrm{\Phi }+\alpha _2\mathrm{\Phi }^2+\mathrm{},`$
$`d`$ $`=`$ $`\mathrm{\Phi }+\beta _0T^2+\beta _1T\mathrm{\Phi }+\beta _2\mathrm{\Phi }^2+\mathrm{},`$ (3.32)
where the dots indicate terms of higher order in the sigma model fields. We found no need for a $`T^2`$ term in the redefinition of tachyon field, such a term would change the cubic and quartic self-couplings of the tachyon in $`V(T)`$. Since $`d`$ gives rise to pure tachyon terms that are quadratic or higher, only at quintic and higher order in $`T`$ will $`V(T)`$ differ from the potential obtained by replacing $`tT`$ in the first line of (3.30). We thus expect that after the field redefinition (3.30) becomes
$$\kappa ^2V=e^{2\mathrm{\Phi }}\left(T^2+1.6018T^33.4693T^4+\mathrm{}\right),$$
(3.33)
at least to quartic order in the fields. We now plug the substitutions (3.3) into the potential (3.30) and compare with (3.33). A number of conditions emerge.
* In order to get the requisite $`T^2\mathrm{\Phi }`$ term we need $`\alpha _1=1`$.
* In order to have a vanishing $`T\mathrm{\Phi }^2`$ term $`\alpha _2=\frac{1}{2}a_{1,2}`$ must be half the coefficient of $`td^2`$ in (3.30).
* Getting the correct $`T^3\mathrm{\Phi }`$ coupling then fixes $`\beta _0=(a_{3,0}a_{3,1})/(2a_{1,2})`$.
* Getting the correct value of $`T^2\mathrm{\Phi }^2`$ fixes $`\beta _1=(1+\frac{3}{2}a_{3,0}a_{1,2}+a_{2,2})/(2a_{1,2})`$. The vanishing of $`T\mathrm{\Phi }^3`$ fixes $`\beta _2=a_{1,3}/(2a_{1,2})`$. All coefficients in (3.3) are now fixed.
* The coefficient of $`\mathrm{\Phi }^4`$, which should be zero, turns out to be $`(a_{0,4}+\frac{1}{4}a_{1,2}^2)0.0389`$, which is small, but does not vanish.
Our inability to adjust the coefficient of $`\mathrm{\Phi }^4`$ was to be expected. The potential (3.30) contains the terms $`t^2+a_{1,2}td^2+a_{0,4}d^4`$ and, to this order, integrating out the tachyon gives an effective dilaton quartic term of $`(a_{0,4}+\frac{1}{4}a_{1,2}^2)`$. With the contribution of the massive fields beyond level four this coefficient in the dilaton effective potential would vanish. This is, in fact, the statement that was verified in . It follows that we need not worry that the quartic term in $`\mathrm{\Phi }`$ do not vanish exactly. Following the steps detailed before we find
$`t`$ $`=`$ $`TT\mathrm{\Phi }\mathrm{\hspace{0.17em}0.4219}\mathrm{\Phi }^2+\mathrm{},`$
$`d`$ $`=`$ $`\mathrm{\Phi }+1.3453T^2+1.1180T\mathrm{\Phi }\mathrm{\hspace{0.17em}0.5646}\mathrm{\Phi }^2+\mathrm{}.`$ (3.34)
In string field theory the dilaton deformation is represented in the $`(d,t)`$ plane by the curve $`(d,t(d))`$, where $`t(d)`$ is the expectation value of the tachyon when the dilaton is set equal to $`d`$. This curve, calculated using the action (3.30), is shown as a solid line in Figure 2. On the other hand, it is clear that $`\mathrm{\Phi }`$ (with $`T=0`$) defines the marginal direction in the effective field theory. Setting $`T=0`$ in (3.3) we find the pair $`(d(\mathrm{\Phi }),t(\mathrm{\Phi }))`$, which must be a parameterization of the flat direction in terms of $`\mathrm{\Phi }`$. This curve is shown as a dashed line in Figure 2. It is a good consistency check that these two curves agree well with each other over a significant fraction of the plot.
### 3.4 Dilaton deformations
In Ref. we computed the effective dilaton potential that arises when we integrate out the tachyon from a potential that includes only quadratic and cubic terms. We found that the domain of definition of this potential is the full real $`d`$ line. This happens because the (marginal) branch $`t(d)`$ that gives the expectation value of $`t`$ for a given value of $`d`$ is well defined for all values of $`d`$. In this section we extend this computation by including higher level fields and higher order interactions. As we will demonstrate, it appears plausible that the domain of definition for the effective dilaton potential remains $`d(\mathrm{},\mathrm{})`$.
The marginal branch is easily identified for small values of the dilaton: as the dilaton expectation value goes to zero all expectation values go to zero. For large enough values of the dilaton the marginal branch may cease to exist, or it may meet another solution branch. If so, we obtain limits on the value of $`d`$. Since the dilaton effective potential is supposed to be flat in the limit of high level, we propose the following criterion. If we encounter a limit value of $`d`$, this value is deemed reliable only if the dilaton potential at this point is not very large. A large value for the potential indicates that the calculation is not reliable because the same terms that are needed to make the potential small could well affect the limit value. In open string field theory a reliable limit value was obtained for the Wilson line parameter: at the limit point the potential energy density was a relatively small fraction of the D-brane energy density. The purely cubic potential for $`t`$ gives a critical point with $`\kappa ^2V0.05774`$. We define $`(d)\frac{|\kappa ^2V(d)|}{0.05774}`$, where $`V(d)`$ is the effective dilaton potential. A critical value of $`d`$ for which $`>1`$ will be considered unreliable.
We start with cubic potentials and then include the elementary quartic interactions level by level. With cubic potentials, the effective dilaton potential is invariant under $`dd`$. With $`𝕍_4^{(3)}`$ dilaton deformations can be arbitrarily large . We then find
* The dilaton potential derived from $`𝕍_8^{(3)}`$ is defined for $`|d|624`$. This is plausible since, at this level, the equations of motion for the level-four fields are linear.
* The dilaton potential derived from $`𝕍_{12}^{(3)}`$ is defined for $`|d|1.71`$. Since $`(\pm 1.71)=42.4`$, there is no reliable limit value.
* The dilaton potential derived from $`𝕍_0^{(4)}`$ is defined for $`|d|4.67`$, where $`(\pm 4.67)=49.5`$. The large value of $``$ indicates that there is no evidence of a limit value.
* The dilaton potential derived from $`𝕍_2^{(4)}`$ is not invariant under $`dd`$. We find a range $`d(\mathrm{},3.124)`$ . Although $`(3.124)=0.387`$, the potential has a maximum with $`=3.325`$ at $`d=1.92`$. This fact makes the limit point $`d=3.124`$ unreliable.
* The dilaton potential derived from $`𝕍_4^{(4)}`$, the highest level potential we have computed fully, is regular for $`d(2.643,\mathrm{\hspace{0.17em}6.415})`$. Since $`(6.415)=1502.4`$ and $`(2.643)=89.2`$, there is no branch cut in the reliable region.
We have also computed the higher level quartic interactions $`td^3`$ and $`d^4`$. We have checked that $`𝕍_4^{(4)}`$, supplemented by those interactions does not lead to branch cuts in the potential for the dilaton. This result, however, is not conclusive. Additional interactions must be included at level six (the level of $`td^3`$) and at level eight (the level of $`d^4`$).
We tested in that cubic and quartic interactions combine to give a vanishing quartic term in the dilaton effective potential. We can ask if the potential for the dilaton becomes flatter as the level of the calculation is increased. We find that it roughly does, but the major changes in the potential are due to the elementary quartic term in the dilaton. For the cubic vertex, the interactions of the type $`d^2M`$, with $`M`$ massive give rise to terms quartic on the dilaton. Other cubic couplings that do not involve the dilaton typically induce $`d^6`$ (and higher order) terms, which play a secondary role in flattening the potential if the quartic terms have not cancelled completely. Therefore, the potentials that arise from $`𝕍_8^{(3)}`$, $`𝕍_{10}^{(3)}`$ and $`𝕍_{12}^{(3)}`$ (without the contribution from level six massive fields) have no obvious difference. The potentials obtained at various levels are shown in Figure 3. The dashed line arises from $`𝕍_4^{(3)}`$, the solid line arises from $`𝕍_8^{(3)}`$, and the thick line arises from $`𝕍_8^{(4)}`$.
## 4 Conclusions
In this paper we have presented some calculations that suggest the existence of a tachyon vacuum for the bulk closed string tachyon of bosonic string theory. We have discussed the physical interpretation using the effective field theory both to suggest the value of the action density at the critical point (zero!) and to obtain rolling solutions that seem consistent with the interpretation of the tachyon vacuum as a state in which there are no closed string states.
The numerical evidence presented is still far from conclusive. A critical point seems to exist and appears to be robust, but it is not all that clear what will happen when the accuracy of the computation is increased. If the action density at the critical point goes to zero it may indeed define a new and nontrivial tachyon vacuum. Conceivably, however, the critical point could approach the perturbative vacuum, in which case there would be no evidence for a new vacuum. Alternatively, if the action density at the critical point remains finite, we would have no interpretation for the result.
Let us consider some additional indirect arguments that support the existence of a closed string tachyon vacuum. The first one arises from the existence of sub-critical bosonic string theories. The evidence in string theory is that most string theories are related by compactifications and/or deformations. It seems very likely that non-critical string theories are also related to critical string theory. It should then be possible to obtain a non-critical string theory as a solution of critical string theory. Certainly the view that $`D=2`$ bosonic string theory is a ground state of the bosonic string has been held as likely . In non-critical string theory the number of space dimensions is reduced (at the expense of a linear dilaton background). The analogy with lower-dimensional D-branes in open string theory seems apt: the branes are solitons of the open string field theory tachyon in which far away from the branes the tachyon sits at the vacuum. It seems plausible that non-critical string theories are solitonic solutions of the closed string theory tachyon. As sketched in Figure 4, far away along the coordinates transverse to the non-critical world-volume, the background would approach the closed string tachyon vacuum. The universality of the tachyon vacuum would imply that a noncritical string theory could be further reduced using the same background configuration used to reduce the original critical theory.
In fact, in the $`p`$-adic open/closed string theory lump solutions of the closed string sector appear to describe spacetimes of lower dimensionality, as explained by Moeller and Schnabl . Indeed, far away from the lump the open string tachyon must be at its vacuum and therefore there are no D-brane solutions with more space dimensions than those of the lump. Away from the lump the closed string tachyon is at its vacuum, and no linearized solutions of the equations of motion exist.
A suggestive argument for zero action at the tachyon vacuum follows from the sigma model approach. As discussed by Tseytlin , it seems likely that the closed string effective action for the spacetime background fields may be written in terms of the partition function $`Z`$ of the two-dimensional sigma model as well as derivatives thereof (this does work for open strings ). The conventional coupling of the world-sheet area to the tachyon $`T`$ results in a partition function and an effective action with a prefactor of $`e^T`$. Thus one expects a tachyon potential of the form $`e^Tg(T)`$ where $`g`$ is a polynomial that begins with a negative quadratic term<sup>4</sup><sup>4</sup>4In , a tachyon potential of the form $`T^2e^T`$ is considered. Complications in fixing the kinetic terms made it unclear if $`T=\mathrm{}`$ was a point in the configuration space (see the discussion below eqn. (4.13)) of . For additional comments on the possible form of the tachyon potential, see Andreev .. In this case, for a tachyon vacuum at $`T\mathrm{}`$ the action goes to zero.
The computations and the discussion presented in this paper have led to a set of testable conjectures concerning the vacuum of the bulk closed string tachyon of bosonic string theory. It seems likely that additional computations, using both string field theory, effective field theory, and conformal field theory will help test these ideas in the near future.
Acknowledgements We are grateful to M. Headrick and A. Sen for many instructive discussions. We would also like to acknowledge useful conversations with K. Hashimoto, H. Liu, N. Moeller, Y. Okawa, M. Schnabl, and A. Tseytlin.
## Appendix A Quartic Computations
### A.1 The setup
We normalize correlators using $`0|c_1\overline{c}_1c_0^{}c_0^+c_1\overline{c}_1|0=1`$ with $`c_0^\pm =\frac{1}{2}(c_0\pm \overline{c}_0)`$. All states in this paper have zero momentum. For convenience, all spacetime coordinates have been compactified and the volume of spacetime is equal to one. To use results from open string field theory, we note that
$$c(z_1)c(z_2)c(z_3)\overline{c}(\overline{w}_1)\overline{c}(\overline{w}_2)\overline{c}(\overline{w}_3)=2c(z_1)c(z_2)c(z_3)_o\overline{c}(\overline{w}_1)\overline{c}(\overline{w}_2)\overline{c}(\overline{w}_3)_o,$$
(A.1)
since open string field theory uses $`c(z_1)c(z_2)c(z_3)_o=(z_1z_2)(z_1z_3)(z_2z_3)`$. Then:
$$c_1\overline{c}_1,c_1\overline{c}_1,c_1\overline{c}_1=2c_1,c_1,c_1_o\overline{c}_1,\overline{c}_1,\overline{c}_1_o=2^3^3=2^6,$$
(A.2)
where $`1/\rho =3\sqrt{3}/41.2990,`$ and $`\rho `$ is the mapping radius of the disks in the three-string vertex.
To construct four-string amplitudes we use antighost insertions
$$=\underset{I=1}{\overset{4}{}}\underset{m=1}{\overset{\mathrm{}}{}}(B_m^Ib_m^J+\overline{C_m^I}\overline{b}_m^I),^{}=\underset{I=1}{\overset{4}{}}\underset{m=1}{\overset{\mathrm{}}{}}(C_m^Ib_m^I+\overline{B_m^I}\overline{b}_m^I),$$
(A.3)
where $`^{}`$ is the $``$-conjugate of $``$. The multilinear function in string field theory is
$$\{\mathrm{\Psi }_1,\mathrm{\Psi }_2,\mathrm{\Psi }_3,\mathrm{\Psi }_4\}\frac{1}{\pi }_{𝒱_{0,4}}𝑑xdy\mathrm{\Sigma }|^{}|\mathrm{\Psi }_1|\mathrm{\Psi }_2|\mathrm{\Psi }_3|\mathrm{\Psi }_4.$$
(A.4)
The first, second, third, and fourth states are inserted at $`0,1,\xi =x+iy,`$ and $`\mathrm{}`$, respectively. Operationally, the fourth state is inserted at $`t=0`$ with $`z=1/t`$, where $`z`$ is the global uniformizer. For further details and explanations the reader should consult . We record that
$$\begin{array}{cc}\hfill B_1^J& =\delta _{3J}/\rho _3,C_1^J=0,\hfill \\ \hfill B_1^I& =\rho _I\beta _I+\frac{1}{2}\rho _3\epsilon _3\delta _{I3},C_1^I=\rho _I\overline{}\beta _I,\hfill \\ \hfill B_2^I& =\frac{1}{6}\rho _I^2(2\beta _I^2\epsilon _I)+\rho _I^2(4\delta _I2\epsilon _I\beta _I+8\beta _I^3)\delta _{3I},C_2^I=\frac{1}{6}\rho _I^2\overline{}(2\beta _I^2\epsilon _I).\hfill \end{array}$$
(A.5)
Here $`\overline{}\frac{}{\overline{\xi }}`$ and $`\frac{}{\xi }`$. Since our string fields are annihilated both by $`b_0`$ and $`\overline{b}_0`$, the coefficients $`B_0^I`$ and $`C_0^I`$ are not needed. Taking note of the vanishing coefficients, we see that for states in the Siegel gauge the antighost factor $``$ is given by
$$=B_1^3b_1^{(3)}+\underset{I=1}{\overset{4}{}}(B_1^Ib_1^J+\overline{C_1^I}\overline{b}_1^I)+\underset{I=1}{\overset{4}{}}(B_2^Ib_2^J+\overline{C_2^I}\overline{b}_2^I)+\mathrm{}.$$
(A.6)
The Strebel quadratic differential on the surfaces determines:
$$\beta _1=\frac{a}{2\xi }\frac{1}{\xi }1,\beta _2=\frac{a2\xi }{2(1\xi )},\beta _3=\frac{a2}{2\xi (\xi 1)},\beta _4=\frac{a}{2}1\xi .$$
(A.7)
Here $`a(\xi ,\overline{\xi })`$ is a function that determines the quadratic differential completely. We also have
$$\begin{array}{cc}\hfill \epsilon _1& =2+\frac{1}{\xi }(a2)+\frac{1}{\xi ^2}\left(2+a\frac{5}{8}a^2\right),\hfill \\ \hfill \epsilon _2& =\frac{5a^2+16\xi (\xi 3)+8a(\xi +3)}{8(\xi 1)^2},\hfill \\ \hfill \epsilon _3& =\frac{16+8a5a^2+24(a2)\xi }{8\xi ^2(\xi 1)^2},\hfill \\ \hfill \epsilon _4& =2+a\frac{5}{8}a^22\xi +a\xi +2\xi ^2.\hfill \end{array}$$
(A.8)
The function $`a(\xi )`$ is known numerically to high accuracy for $`\xi 𝒜`$, where $`𝒜`$ is a specific subspace of $`𝒱_{0,4}`$ described in detail in Figures 3 and 6 of ref. . The full space $`𝒱_{0,4}`$ is obtained by acting on $`𝒜`$ with the transformations generated by $`\xi 1\xi `$ and $`\xi 1/\xi `$, together with complex conjugation $`\xi \overline{\xi }`$. In fact $`𝒱_{0,4}`$ contains twelve copies of $`𝒜`$. Let $`f(𝒜)`$ denote the region obtained by mapping each point $`\xi 𝒜`$ to $`f(\xi )`$. Then $`𝒱_{0,4}`$ is composed of the six regions
$$𝒜,\frac{1}{𝒜},1𝒜,\frac{1}{1𝒜},1\frac{1}{𝒜},\frac{𝒜}{1𝒜},$$
(A.9)
together with their complex conjugates. The values of $`a`$ in these regions follow from the values of $`a`$ on $`𝒜`$ via the relations
$$a(1\xi )=4a(\xi ),a\left(\frac{1}{\xi }\right)=\frac{a(\xi )}{\xi },a(\overline{\xi })=\overline{a(\xi )}.$$
(A.10)
For states of the form $`|M_i=𝒪_ic_1\overline{c}_1|0`$, where $`𝒪_i`$ is built with matter oscillator, one finds
$$\{M_1,M_2,M_3,M_4\}=\frac{2}{\pi }_{𝒱_{0,4}}\frac{dxdy}{(\rho _1\rho _2\rho _3\rho _4)^2}𝒪_1𝒪_2𝒪_3𝒪_4_\xi .$$
(A.11)
Here $`𝒪_1𝒪_2𝒪_3𝒪_4_\xi h_1𝒪_1h_2𝒪_2h_3𝒪_3h_4𝒪_4_{\mathrm{\Sigma }_\xi },`$ where the right-hand side is a matter correlator computed after the local operators $`𝒪_i`$ have been mapped to the uniformizer.
### A.2 Couplings of dilatons and tachyons
Elementary contribution to $`t^3d`$. We insert the dilaton on the moving puncture to make the integration identical over each of the 12 regions of the moduli space. Since all the states inserted on the fixed punctures have ghost oscillators $`c_1\overline{c}_1`$, the antighost factor $`^{}`$ is only supported on the moving puncture:
$`^{}(c_1c_1\overline{c}_1\overline{c}_1)^{(3)}|0=(B_1^3C_1^3+\overline{B_1^3}\overline{C_1^3})|0=(\overline{}\beta _3+\overline{\beta }_3)|0.`$ (A.12)
There are no matter operators, thus the correlator just involves the ghosts:
$`\mathrm{\Sigma }_P|^{}|T|T|D|T`$ $`=`$ $`(\overline{}\beta _3+\overline{\beta }_3)(c_1\overline{c}_1)^{(1)}(c_1\overline{c}_1)^{(2)}(c_1\overline{c}_1)^{(4)}`$ (A.13)
$`=`$ $`(\overline{}\beta _3+\overline{\beta }_3){\displaystyle \frac{2}{(\rho _1\rho _2\rho _4)^2}}.`$
Using (A.4), the amplitude is:
$$\{T^3D\}=\frac{24}{\pi }_𝒜𝑑x𝑑y(\overline{}\beta _3+\overline{\beta }_3)\frac{1}{(\rho _1\rho _2\rho _4)^2}=\mathrm{\hspace{0.17em}23.2323}.$$
(A.14)
The contribution to the potential is $`\kappa ^2V=\frac{4}{4!}\{T^3D\}t^3d=3.8721t^3d.`$
Elementary contribution to $`t^2d^2`$. We insert the dilatons at $`z_2=1`$ and $`z_3=\xi `$. The amplitude to be integrated is identical to the ghost part of the amplitude for the quartic interaction $`a^2d^2`$, as given in , equation (4.9):
$$\mathrm{\Sigma }|^{}|T|D|D|T=\frac{2}{(\rho _1\rho _4)^2}(\overline{}\beta _2(\overline{\xi }\overline{\beta }_3)\beta _2\overline{}(\overline{\xi }\overline{\beta }_3)+\text{-conj}).$$
(A.15)
The four-point amplitude is then
$`\{T^2D^2\}={\displaystyle \frac{4}{\pi }}{\displaystyle _{𝒱_{0,4}}}{\displaystyle \frac{dxdy}{(\rho _1\rho _4)^2}}\text{Re}\left(\overline{}\beta _2(\overline{\xi }\overline{\beta }_3)\beta _2\overline{}(\overline{\xi }\overline{\beta }_3)\right).`$ (A.16)
Since we have the same states on punctures one and four, and these punctures are exchanged by the transformation $`z1/z`$, the integral over $`𝒜`$ gives the same contribution as the integral over $`1/𝒜`$. The conjugation properties of the amplitude also imply that $`\overline{𝒜}`$ contributes the same as $`𝒜`$. Consequently, the four regions $`𝒜,\mathrm{\hspace{0.17em}1}/𝒜,\overline{𝒜},`$ and $`1/\overline{𝒜}`$ all give the same contribution. To get the full amplitude we must multiply the contributions of $`𝒜`$, of $`1𝒜`$, and $`11/𝒜`$ by four:
$`\{T^2D^2\}=4{\displaystyle \frac{4}{\pi }}\left[{\displaystyle _𝒜}+{\displaystyle _{1𝒜}}+{\displaystyle _{11/𝒜}}\right]{\displaystyle \frac{dxdy}{(\rho _1\rho _4)^2}}\text{Re}\left(\overline{}\beta _2(\overline{\xi }\overline{\beta }_3)\beta _2\overline{}(\overline{\xi }\overline{\beta }_3)\right).`$ (A.17)
The transformation laws given in Appendix B of allow one to rewrite the second and third integrals as integrals over $`𝒜`$, where they can be easily evaluated. We find
$`\{T^2D^2\}=4(0.2410+0.4031+1.2065)=5.4726.`$ (A.18)
The contribution to potential is $`\kappa ^2V=\frac{6}{4!}\{T^2D^2\}t^2d^2=1.3682t^2d^2.`$
Elementary contribution to $`td^3`$. The tachyon field is inserted at $`z_3=\xi `$. We then have
$`^{}(c_1\overline{c}_1)^{(3)}D^{(1)}D^{(2)}D^{(4)}|0`$ (A.19)
$`=`$ $`\left\{B_1^3b_1^{(3)}+{\displaystyle \underset{J3}{}}\left(B_1^Jb_1^{(J)}+\overline{C_1^J}\overline{b}_1^{(J)}\right)\right\}\left\{\overline{B_1^3}\overline{b}_1^{(3)}+{\displaystyle \underset{J3}{}}\left(\overline{B_1^J}\overline{b}_1^{(J)}+C_1^Jb_1^{(J)}\right)\right\}(c_1\overline{c}_1)^{(3)}D^{(1)}D^{(2)}D^{(4)}|0`$
$`=`$ $`{\displaystyle \underset{IJK3}{}}(\frac{1}{2}B_1^3C_1^ID^{(J)}D^{(K)}c_1^{(I)}\overline{c}_1^{(3)}+B_1^IC_1^J(c_1\overline{c}_1)^{(3)}c_1^{(I)}c_1^{(J)}(\overline{c}_1\overline{c}_1)^{(K)})|0\text{-conj}.`$
Therefore, the correlator $`𝒞_{td^3}=\mathrm{\Sigma }|^{}TD^3|0`$ is:
$$𝒞_{td^3}=\underset{IJK3}{}B_1^3C_1^I(\overline{c}_1\overline{c}_1)^{(J)}(c_1c_1)^{(K)}c_1^{(I)}\overline{c}_1^{(3)}+B_1^IC_1^J(\overline{c}_1\overline{c}_1)^{(K)}c_1^{(I)}c_1^{(J)}(c_1\overline{c}_1)^{(3)}+\text{-conj}.$$
Factorizing into holomorphic and antiholomorphic parts we get
$$𝒞_{td^3}=2\underset{IJK3}{}(B_1^3C_1^IB_{KI}(B_{J3})^{}B_1^IC_1^JD_{IJ}(B_{K3})^{})+\text{-conj},$$
(A.20)
where $`B_{IJ}(c_1c_1)^{(I)},c_1^{(J)}`$ was introduced and evaluated in , eqns. (4.18), (4.20), and (4.21). Additionally,
$$D_{IJ}c_1^{(I)},c_1^{(J)},c_1^{(3)}=\frac{z_{IJ}z_{I3}z_{J3}}{\rho _I\rho _J\rho _3},D_{I4}=D_{4I}=\frac{z_{I3}}{\rho _I\rho _3\rho _4},I,J4.$$
(A.21)
The full amplitude is
$$\{TD^3\}=\frac{12}{\pi }_𝒜𝑑x𝑑y𝒞_{td^3}=5.7168.$$
(A.22)
The contribution to the potential is $`\kappa ^2V=\frac{4}{4!}\{TD^3\}td^3=0.9528td^3.`$
### A.3 Couplings of tachyon to massive fields
In all cases the massive field will be inserted on the moving puncture $`z_3=\xi `$.
Elementary contribution to $`t^3f_1`$. With $`F_1c_1\overline{c}_1`$ inserted at $`z_3=\xi `$ we find:
$$^{}(c_1\overline{c}_1)^{(3)}|0=(C_1^3\overline{C_1^3}B_1^3\overline{B_1^3})|0.$$
(A.23)
$$\{T^3F_1\}=\frac{12}{\pi }_𝒜𝑑x𝑑y\frac{2}{(\rho _1\rho _2\rho _4)^2}(C_1^3\overline{C_1^3}B_1^3\overline{B_1^3})=2.6261.$$
(A.24)
The contribution to the potential is: $`\kappa ^2V=\frac{4}{4!}\{T^3F_1\}t^3f_1=0.4377t^3f_1.`$
Elementary contribution to $`t^3f_2`$. With $`F_2c_1\overline{c}_1L_2\overline{L}_2`$ at $`z_3=\xi `$, the ghost part is that of the four-tachyon amplitude (eqn. (3.34) of ). With $`w=0`$ corresponding to $`z=z_3`$, and $`S(z,w)`$ denoting the Schwarzian derivative, the holomorphic matter correlator is:
$$\begin{array}{cc}\hfill L_2^{(3)}& =T^{(3)}(w=0)=\rho _3^2T(z_3)+\frac{26}{12}S(z,w)=\frac{13}{6}\rho _3^2(2\beta _3^2\epsilon _3).\hfill \end{array}$$
(A.25)
Therefore, the amplitude is
$$\{T^3F_2\}=\frac{24}{\pi }_𝒜\frac{dxdy}{(\rho _1\rho _2\rho _3\rho _4)^2}\left|\frac{13}{6}\rho _3^2(2\beta _3^2\epsilon _3)\right|^2=\mathrm{\hspace{0.17em}337.571}.$$
(A.26)
The contribution to the potential is $`\kappa ^2V=\frac{4}{4!}\{T^3F_2\}t^3f_2=56.262t^3f_2.`$
Elementary contribution to $`t^3f_3`$. With $`L_2c_1\overline{c}_1`$ inserted at $`z_3=\xi `$ we find
$$^{}(c_1\overline{c}_1)^{(3)}|0=B_1^3\overline{B_1^3}|0.$$
(A.27)
$$𝒞_{t^3f_3}\mathrm{\Sigma }|^{}TT(c_1\overline{c}_1)^{(3)}T|0L_2^{(3)}=\frac{2B_1^3\overline{B_1^3}}{(\rho _1\rho _2\rho _4)^2}\frac{13}{6}\rho _3^2(2\beta _3^2\epsilon _3).$$
(A.28)
With $`F_3L_2c_1\overline{c}_1+c_1\overline{L}_2\overline{c}_1`$, the string amplitude relevant to $`t^3f_3`$ is:
$$\{T^3F_3\}=\frac{12}{\pi }_𝒜𝑑x𝑑y(𝒞_{t^3f_3}+𝒞_{t^3f_3}^{})=78.1432.$$
(A.29)
The contribution to the potential is: $`\kappa ^2V=\frac{4}{4!}\{T^3F_3\}t^3f_3=13.024t^3f_3.`$
Elementary contribution to $`t^3g_1`$. With $`b_2c_1\overline{c}_2\overline{c}_1`$ at $`z_3=\xi `$, one finds
$$^{}(b_2c_1\overline{c}_2\overline{c}_1)^{(3)}|0=\overline{C_2^3}\overline{B_1^3}(c_1b_2)^{(3)}|0.$$
(A.30)
The state $`c_1b_2|0`$ is created by the non-primary ghost current $`j(z)=cb(z)`$ by acting on the vacuum. For the ghost current
$$j(w)=j(z)\frac{dz}{dw}\frac{3}{2}\frac{z^{\prime \prime }}{z^{}}j(w=0)=\rho _3(j(z_3)3\beta _3).$$
(A.31)
We thus have the correlator:
$$\begin{array}{cc}\hfill 𝒞_{t^3g_1}& \mathrm{\Sigma }|^{}TT(b_2c_1\overline{c}_2\overline{c}_1)^{(3)}T|0\hfill \\ \hfill =& \overline{C_2^3}\overline{B_1^3}\frac{1}{(\rho _1\rho _2\rho _4)^2}c\overline{c}(0)c\overline{c}(1)\rho _3(j(z_3)3\beta _3)c\overline{c}(t=0)\hfill \\ \hfill =& \overline{C_2^3}\overline{B_1^3}\frac{\rho _3}{(\rho _1\rho _2\rho _4)^2}2\left(\frac{1}{\xi }+\frac{1}{\xi 1}3\beta _3\right).\hfill \end{array}$$
(A.32)
With $`G_1b_2c_1\overline{c}_2\overline{c}_1c_2c_1\overline{b}_2\overline{c}_1`$, the amplitude relevant for the $`t^3g_1`$ coupling is
$$\{T^3G_1\}=\frac{12}{\pi }_𝒜𝑑x𝑑y(𝒞_{t^3g_1}+𝒞_{t^3g_1}^{})=1.6350.$$
(A.33)
The contribution to the potential is $`\kappa ^2V=\frac{4}{4!}\{T^3G_1\}t^3g_1=0.2725t^3g_1`$. |
warning/0506/hep-th0506194.html | ar5iv | text | # Observing an open FRW de Sitter universe living in a Minkowski spacetime.
## I Introduction
The WMAP results WMAP combined with earlier cosmological observations shows that we are living in an accelerating universe. The currently observed lumpiness in the temperature of the cosmic microwave background is just right for a flat universe though there are also some evidences that our universe is spatially open Gott . The great simplifying fact of cosmology is that the universe appears to be homogeneous and isotropic along a preferred set of spatial hypersurfaces 16 . Of course homogeneity and isotropy are only approximate, but they become increasingly good approximations on larger length scales, allowing us to describe spacetime on cosmological scales by the Robertson-Walker metric. Constructing four dimensional de Sitter vacuum as a string theory (M-theory) solution has been a long standing challenge. An outstanding example of string theory models of de Sitter vacua are the KKLT models KKLT with an exponentially large number of stable and metastable vacua without supersymmetry or with $`𝒩=1`$ supersymmetry in four dimensions, the landscape Sus . In KKLT models, metastable de Sitter vacua of type IIB string theory is constructed by adding $`\overline{\text{D3}}`$-branes to the GKP GKP model of highly warped IIB compactifications with nontrivial NS and RR three-form fluxes after certain fine tuning of the fluxes.
Recently, we realized that it is possible to observe a de Sitter universe while living in a flat background. This proposal was the consequence of a simple observation: the fluctuations of the scalar field around the classical trajectory of an unstable massless $`\varphi ^4`$ model in four dimensional flat Euclidean spacetime is governed by a conformally coupled scalar field theory in four dimensional de Sitter background Solitons . This classical trajectory is the Fubini vacua of the classically conformal-invariant scalar field theories. In Fubini S. Fubini verified that critical scalar theories possess a classical vacua with $`O(D,1)`$ symmetry in which the expectation value of the scalar field is non-vanishing. The motivation to study such a classical vacua at that time was ”to introduce a fundamental scale of hadron phenomena by means of dilatation non-invariant vacuum state in the frame work of a scale invariant Lagrangian field theory” Fubini . This result is interesting due to its uniqueness. In four dimensions, in principle, one can consider two classes of critical (classically scale-free ) scalar field theories i.e. massless $`\varphi ^4`$ models on Euclidean spacetime with $`g`$, the coupling constant, either positive or negative (we assume the potential $`V(\varphi )=\frac{g}{4}\varphi ^4`$). Although scalar theory with $`g>0`$ seems to be not physical as the potential is not bounded from below but in this case, the Euler-Lagrange equation of motion has an interesting classical solution say $`\varphi _0`$ with finite action $`S[\varphi _0]g^1`$. For $`g<0`$, one can still consider a solution like $`\varphi _0`$ obtained by an analytic continuation from $`g>0`$ to $`g<0`$ region. But such a solution is singular on the surface of a sphere which radius is proportional to $`g`$. Consequently the action $`S[\varphi _0]`$ is infinite and $`\varphi _0`$ can not be considered as a classical trajectory. For $`g>0`$ it is shown in Solitons that the information geometry of the moduli space of $`\varphi _0`$ given by Hitchin formula Hit is Euclidean $`\text{AdS}_5`$,
$$𝒢_{IJ}d\theta ^Id\theta ^J=\frac{1}{\beta ^2}\left(d\beta ^2+da^2\right),$$
(1)
where $`\theta ^I\{\beta ,a^\mu \}`$ and $`I=1,\mathrm{},5`$. The moduli here are $`a_\mu `$’s the location of the center of $`\varphi _0`$ and $`\beta `$ which is proportional to the inverse of the size of $`\varphi _0`$. This resembles the information geometry of SU(2) instantons. In addition $`V(\varphi _0)`$ can be shown to be proportional to the SU(2) one-instanton density Blau . Interestingly, $`\varphi _0`$ is the bulk to boundary propagator in the $`\text{AdS}_5`$ geometry of the moduli space.
In ref. U1 we generalized the model to scalar theory coupled to U(1) gauge field. Such a generalization is essential as it shows how by optical observations people living in a flat Euclidean space observe a de Sitter geometry for their universe. In this model the massless scalar field is charged though we have not observed light charged scalars. This problem can be resolved noting that as we will show, in this model observations are made in a de Sitter background in which the scalar field appears to be conformally coupled to the de Sitter background. Therefore its mass is proportional to the scalar curvature $`R`$ of the observed universe. Using the WKB approximation, the lifetime of the observed de Sitter background is calculated in U1 and is shown to be proportional to $`e^{g^1}`$. (To my knowledge, this result is given for the first time in a beautiful paper by Coleman where $`\varphi _0`$ is called a ”bounce” Coleman .) Consequently in the weak coupling limit $`g{}_{}{}^{+}0`$ the lifetime increases exponentially.
In this paper we study the critical scalar theory on Minkowski spacetime. Here $`\varphi _0`$ is singular on a hyperbola in the timelike region which asymptotes to the lightcone. The total energy of $`\varphi _0`$, measured by an observer located at the center of $`\varphi _0`$ is conserved and vanishing though the energy density is a function of space and time. The energy density diverges in the neighborhood of singularity causing a gravitational collapse when the scalar theory is coupled to gravity. Fortunately the singularity is safe. On the one hand in Minkowski spacetime, the distance between the observers and the singularity is proportional to $`\beta `$. On the other hand there is some mechanism of $`\beta `$ transition in the model: larger $`\varphi _0`$’s decay to smaller ones due to say thermal fluctuations around $`\varphi _0`$ and finally there remains only a gas of zero sized bubbles which corresponds to $`\beta \mathrm{}`$. The mechanism of such transition is not clear yet but its phenomenology, probably is similar to that of the discretuum of possible de Sitter vacua in KKLT models fall . Therefore, for the most stable $`\varphi _0`$ solution, the singularity is located at infinite future and is out of reach. Furthermore, from the observers point of view, the observable universe is an open de Sitter space which horizon is located on the singularity. Therefore they do not see the singularity at all for any value of $`\beta `$! Of course they should feel some back reactions when the scalar theory is coupled to gravity caused by the $`\beta `$ transition. The energy density $`\rho `$ and pressure $`p`$ that they measure are constants satisfying the dark energy equation of state $`\rho =p=\mathrm{\Lambda }`$ ($`8\pi G=1`$), where $`\mathrm{\Lambda }`$ is the cosmological constant. A question here is the value of $`\mathrm{\Lambda }`$ or equivalently $`R`$, the curvature scalar. At classical level, $`R`$ is not determined in the critical scalar model as is expected. Because the theory is classically scale free. The quantum theory is not scale free due to loop corrections. Therefore quantum corrections are the hopeful candidates to give the value of the observed scalar curvature. The details are not clear for us yet and we postpone it to future works.
In Minkowski spacetime in the case of critical scalar model with negative coupling constant the singularity of $`\varphi _0`$ is a hyperbola in the spacelike region which asymptotes to the lightcone. One can easily verify that in this case the total energy for existence of $`\varphi _0`$ is infinite. Thus, similar to the Euclidean case, one can conclude that $`\varphi _0`$ uniquely exists only in the unstable ($`g>0`$) critical scalar model.
The organization of the paper is as follows. In the next section we study the critical scalar theory on flat Euclidean background and determine the role of the moduli $`\beta `$ in the stability of the solutions. We show that by recasting the scalar theory in terms of new fields $`\stackrel{~}{\varphi }=\varphi \varphi _0`$ at the end of the day one obtains a $`\varphi ^4`$ model conformally coupled to a de Sitter background. In section III we switch to the Minkowski spacetime by a Wick rotation $`tit`$ and study the observed de Sitter universe in terms of the Robertson-Walker metric.
## II Critical scalar theory in $`D=4`$ Euclidean space
In this section we study scalar field theories in four dimensional Euclidean space invariant under rescaling transformation $`xx^{}=\lambda x`$, $`\lambda >0`$. There are in general three scale-free scalar theories: $`\varphi ^4`$ model in $`D=4`$ and $`\varphi ^3`$ and $`\varphi ^6`$ models in $`D=6,3`$ respectively. In this paper we only study $`\varphi ^4`$ model in $`D=4`$ though the main result of this paper can be simply generalized to the other two scalar models. The action in Euclidean space is
$$S[\varphi ]=d^4x\left(\frac{1}{2}\delta ^{\mu \nu }_\mu \varphi _\nu \varphi \frac{g}{4}\varphi ^4\right)$$
(2)
where we assume $`g>0`$. Consequently the potential $`V(p)\varphi ^4`$ and is not bounded from below. The Kronecker delta symbol $`\delta ^{\mu \nu }`$ stands for the metric of flat Euclidean space. The corresponding equation of motion is a non-linear Laplace equation $`^2\varphi +g\varphi ^3=0`$, where $`^2=\delta ^{\mu \nu }_\mu _\nu `$. One can easily show that for $`g>0`$, a solution of the non-linear Laplace equation is
$$\varphi _0(x;\beta ,a^\mu )=\sqrt{\frac{8}{g}}\frac{\beta }{\beta ^2+(xa)^2},$$
(3)
where $`(xa)^2=\delta _{\mu \nu }(xa)^\mu (xa)^\nu .`$ $`\beta `$ and $`a^\mu `$ are undetermined parameters describing the the size and location of $`\varphi _0`$. These moduli are consequences of symmetries of the action i.e. invariance under rescaling and translation. The information geometry of the moduli space, given by Hitchin formula Hit
$$𝒢_{IJ}=\frac{1}{N}d^4x_0_I\left(\mathrm{log}_0\right)_J\left(\mathrm{log}_0\right),$$
(4)
is an Euclidean $`\text{AdS}_5`$ space (1). $`N=\frac{4^3}{5}d^4x_0`$ is a normalization constant and $`_0=\frac{g}{4}\varphi _0^4`$ is the Lagrangian density calculated at $`\varphi =\varphi _0`$. The moduli $`a^\mu `$ are present since the action is invariant under translation. The existence of $`\beta `$ is the result of invariance under rescaling U1 .
To my knowledge, the solution $`\varphi _0`$ is obtained for the first time by Fubini. He looked for a solution of the equation of motion ”in which the vacuum expectation value of the field $`\varphi (x)`$ is non-vanishing” Fubini . He verified that this vacua is not invariant under the Poincare group but is invariant under the de Sitter group $`O(3,1)`$. Consequently by recasting the action in terms of new fields $`\stackrel{~}{\varphi }=\varphi \varphi _0`$ one expects to obtain, after some field redefinitions, a scalar theory in de Sitter background. In fact the action in terms of $`\stackrel{~}{\varphi }`$ is,
$$S[\varphi ]=S[\varphi _0]+S_{\text{free}}[\stackrel{~}{\varphi }]+S_{\text{int}}[\stackrel{~}{\varphi }],$$
(5)
where $`S[\varphi _0]=d^4x_0=\frac{8\pi ^2}{3g}`$, and
$$S_{\text{free}}[\stackrel{~}{\varphi }]=d^4x\left(\frac{1}{2}\delta ^{\mu \nu }_\mu \stackrel{~}{\varphi }_\nu \stackrel{~}{\varphi }+\frac{1}{2}M^2(x)\stackrel{~}{\varphi }^2\right)$$
(6)
in which,
$$M^2(x)=3g\varphi _0^2=24\frac{\beta ^2}{\left(\beta ^2+(xa)^2\right)^2}.$$
(7)
These equations show that $`\varphi _0`$ is a metastable local minima of the action. This can also be verified explicitly by numerical analysis of action (2), see ref. U1 . Equation (6) can be used to show that the stability increases as $`\beta \mathrm{}`$. In fact if we calculate the variation of action at the stationary point $`\varphi _0(\beta )`$ for different values of the moduli $`\beta _1`$ and $`\beta _2`$, under variation $`\delta \varphi `$, from Eqs.(6,7) one verifies that,
$`\mathrm{\Delta }S`$ $`=`$ $`\delta S|_{\beta _1}\delta S|_{\beta _2}`$ (8)
$``$ $`{\displaystyle d^4x\left(\varphi _0(\beta _2)^2\varphi _0(\beta _1)^2\right)\delta \varphi ^2}+𝒪(\delta \varphi ^3).`$
For simplicity we assume that $`a_i^\mu =0`$, $`i=1,2`$. Therefore $`\mathrm{\Delta }S`$ is proportional to,
$$(\beta _1^2\beta _2^2)_0^{\mathrm{}}𝑑x\frac{x^3(x^4+\beta _1^2\beta _2^2)}{(\beta _1^2+x^2)^2(\beta _1^2+x^2)^2}\delta \varphi ^2.$$
(9)
For $`\delta \varphi `$ with compact support, i.e. $`\delta \varphi =0`$ if $`\left|x\right|>\sqrt{\beta _1\beta _2}`$ the integral above is positive therefore $`\mathrm{\Delta }S(\beta _1^2\beta _2^2)`$. As far as $`\varphi _0`$ is a metastable local minima there exist $`\delta \varphi `$ with compact support such that $`\delta S|_{\beta _i}>0`$ $`i=1,2`$. Consequently if $`\beta _1>\beta _2`$ then $`\delta S|_{\beta _1}>\delta S|_{\beta _2}>0`$. One can convince herself/himself that for some $`\delta \varphi `$ one obtains $`\delta S|_{\beta _2}<0`$ while $`\delta S|_{\beta _1}>0`$. Consequently one concludes that there is a transition $`\beta _2\beta _1`$ induced by say, thermal fluctuations. In addition the stability increases as $`\beta \mathrm{}`$.
The mass term in Eq.(6) can be interpreted as interaction with the background $`\varphi _0`$. Now recall that in general, by inserting $`\stackrel{~}{\varphi }=\sqrt{\mathrm{\Omega }}\overline{\varphi }`$ and $`\delta _{\mu \nu }=\mathrm{\Omega }^1g_{\mu \nu }`$ in the action $`S[\stackrel{~}{\varphi }]=d^4x\frac{1}{2}\delta ^{\mu \nu }_\mu \stackrel{~}{\varphi }_\nu \stackrel{~}{\varphi }`$, one obtains,
$$S[\stackrel{~}{\varphi }]=d^4x\sqrt{g}\left(\frac{1}{2}g^{\mu \nu }_\mu \overline{\varphi }_\nu \overline{\varphi }+\frac{1}{2}\xi R\overline{\varphi }^2\right),$$
(10)
i.e. a scalar theory on conformally flat background given by the metric $`g_{\mu \nu }=\mathrm{\Omega }\delta _{\mu \nu }`$ in which $`\mathrm{\Omega }>0`$ is an arbitrary $`𝒞^{\mathrm{}}`$ function. $`R`$ is the scalar curvature of the background and $`\xi =\frac{1}{6}`$ is the conformal coupling constant. For details see Ted or appendix C of Solitons . Thus, defining $`\overline{\varphi }=\mathrm{\Omega }^{\frac{1}{2}}\stackrel{~}{\varphi }`$, one can show that $`S_{\text{free}}[\stackrel{~}{\varphi }]`$ given in Eq.(6) is the action of the scalar field $`\overline{\varphi }`$ on some conformally flat background,
$$S_{\text{free}}[\stackrel{~}{\varphi }]=d^4x\sqrt{\left|g\right|}\left(\frac{1}{2}g^{\mu \nu }_\mu \overline{\varphi }_\nu \overline{\varphi }+\frac{1}{2}(\xi R+m^2)\overline{\varphi }^2\right).$$
(11)
with metric
$$g_{\mu \nu }=\mathrm{\Omega }\delta _{\mu \nu },\mathrm{\Omega }=\frac{M^2(x)}{m^2},$$
(12)
where $`m^2`$ is the mass of $`\overline{\varphi }`$ (undetermined) and $`M^2(x)`$ is given in Eq.(7). This result is surprising as one can show that the Ricci tensor $`R_{\mu \nu }=\mathrm{\Lambda }g_{\mu \nu }`$, where $`\mathrm{\Lambda }=\frac{m^2}{2}>0`$ as far as $`\mathrm{\Omega }>0`$. Consequently $`\overline{\varphi }`$ lives in a four dimensional de Sitter space which scalar curvature $`R=2m^2`$. The interacting part of the action, $`S_{\text{int}}[\stackrel{~}{\varphi }]=d^4x\sqrt{\left|g_{\mu \nu }\right|}_{\text{int}}`$ is well-defined in terms of $`\overline{\varphi }`$ on the corresponding $`\text{dS}_4`$:
$$_{\text{int}}=g\sqrt{\frac{m^2}{3g}}\overline{\varphi }^3\frac{g}{4}\overline{\varphi }^4.$$
(13)
Interestingly after a shift of the scalar field $`\overline{\varphi }\overline{\varphi }\sqrt{\frac{m^2}{3g}}`$ the action (5) can be written in the $`\text{dS}_4`$ as follows:
$$S[\overline{\varphi }]=d^4x\sqrt{\left|g\right|}\left(\frac{1}{2}g^{\mu \nu }_\mu \overline{\varphi }_\nu \overline{\varphi }+\frac{1}{2}(\xi R)\overline{\varphi }^2\frac{g}{4}\overline{\varphi }^4\right).$$
(14)
This is a scalar theory in a de Sitter background with reversed Mexican hat potential. In a similar way, by recasting the critical scalar theory minimally coupled to $`U(1)`$ gauge field in terms of fluctuations around the classical solution $`\varphi =\varphi _0`$ and $`A_\mu =0`$, one verifies that the action
$$S=d^4x\left(\left|D_\mu \varphi \right|^2\frac{g}{2}\left|\varphi \right|^4\right)+S_A,$$
(15)
is equivalent to
$$S=S[\varphi _0]+S[\overline{\varphi },A_\mu ]+S_A,$$
(16)
where
$`S[\overline{\varphi },A_\mu ]={\displaystyle d^4x\sqrt{g}\left(\frac{1}{2}g^{\mu \nu }D_\mu \overline{\varphi }D_\nu \overline{\varphi }^{}+V(\overline{\varphi })\right)},`$ (17)
$`V(\overline{\varphi })=\frac{1}{2}\xi R\left|\overline{\varphi }\right|^2\frac{g}{4}\left|\overline{\varphi }\right|^4`$ and $`S_A`$ is the Kinetic term for the gauge field,
$`S_A`$ $`=`$ $`{\displaystyle \frac{1}{4}}{\displaystyle d^4xF_{\mu \nu }F_{\rho \sigma }\delta ^{\rho \mu }\delta ^{\sigma \nu }}`$ (18)
$`=`$ $`{\displaystyle \frac{1}{4}}{\displaystyle d^4x\sqrt{g}g^{\mu \rho }g^{\nu \sigma }F_{\mu \nu }F_{\rho \sigma }}.`$
$`F_{\mu \nu }`$ in the first equality above is the field strength in Minkowski spacetime. In the second equality $`F_{\mu \nu }`$ should be understood as the field strength on the de Sitter space U1 . It should be noted that under the conformal transformation $`g_{\mu \nu }\mathrm{\Omega }g_{\mu \nu }`$, in four dimensions $`A_\mu A_\mu `$.
## III The critical scalar theory in Minkowski spacetime
The critical scalar theory in four dimensional Minkowski spacetime is given by the action
$$S[\varphi ]=d^4x\left(\frac{1}{2}\eta ^{\mu \nu }_\mu \varphi _\nu \varphi +\frac{g}{4}\varphi ^4\right),$$
(19)
where $`\eta _{\mu \nu }=(+,,,)`$ and $`g>0`$. The equation of motion is a non-linear wave equation $`\eta ^\mu _\mu _\nu \varphi g\varphi ^3=0`$ which has the solution
$$\varphi _0=\sqrt{\frac{8}{g}}\frac{\beta }{\beta ^2(ta^0)^2+\left|\stackrel{}{x}\stackrel{}{a}\right|^2},$$
(20)
where $`\stackrel{}{x}𝐑^3`$. Here on we assume $`a^\mu =0`$ for simplicity. $`\varphi _0`$ is singular on the hyperbola $`t^2=x^2+\beta ^2`$ and we define its distance to an observer located on the origin to be given by $`\beta `$. The Hamiltonian density $``$ corresponding to $`\varphi _0`$, is,
$$=\frac{16\beta ^2}{g}\frac{t^2+x^2\beta ^2}{(t^2+x^2+\beta ^2)^4}$$
(21)
which tends to infinity in the vicinity of the singularity. As is explained in the introduction, using the results of section II and the arguments after Eq.(9), we now that the most stable $`\varphi _0`$ is the zero-sized one, corresponding to $`\beta \mathrm{}`$. Therefore the singularity is safe when the scalar theory is coupled to gravity. For $`t<\beta `$ one can calculate, say, the total vacuum energy $`H=d^3x`$ corresponding to $`\varphi _0`$ which is surprisingly vanishing, $`H=0`$. Repeating the calculations of section II one verifies that observers located at the origin of the Minkowski spacetime observe a de Sitter space given by the conformally flat metric,
$$ds^2=\frac{12\beta ^2}{\mathrm{\Lambda }}\frac{1}{(\beta ^2t^2+x^2)^2}(dt^2+d\stackrel{}{x}^2),$$
(22)
where $`\mathrm{\Lambda }>0`$ is the cosmological constant. This metric can be obtained using Eq.(12) after a Wick rotation $`tit`$. We use a different set of coordinates in order to describe the observed de Sitter space with FRW metric to see whether it is open, closed or flat. Defining, coordinates $`u`$, $`\rho `$ and $`z_i`$, $`i=1,2,3`$ by the relations $`z_i^2=1`$, $`t=u\mathrm{cosh}\rho `$ and $`x_i=u\mathrm{sinh}\rho z_i`$ useful to describe the timelike region $`t>\left|\stackrel{}{x}\right|`$, one obtains,
$$ds^2=\frac{12\beta ^2}{\mathrm{\Lambda }}\frac{1}{(\beta ^2u^2)^2}\left(du^2+u^2(d\rho ^2+\mathrm{sinh}\rho ^2dz_i^2)\right).$$
(23)
we define a time coordinate $`\tau `$ by the relation $`d\tau =\left(\beta ^2u^2\right)^1du`$. Thus one obtains,
$$\tau =\{\begin{array}{cccc}\frac{1}{\beta }\mathrm{coth}^1\frac{u}{\beta }\hfill & & & u>\beta ,\hfill \\ \frac{1}{\beta }\mathrm{tanh}^1\frac{u}{\beta }\hfill & & & u<\beta ,\hfill \end{array}$$
(24)
and
$$ds^2=\frac{12\beta ^2}{\mathrm{\Lambda }}\left(d\tau ^2+\frac{\mathrm{sinh}^2(2\beta \tau )}{4\beta ^2}(d\rho ^2+\mathrm{sinh}^2\rho dz_i^2)\right)$$
(25)
One can call the region $`u<\beta `$ which can be observed by observers located on the origin the south pole and the $`u>\beta `$ region the north pole, a known terminology in de Sitter geometry. The south pole and north pole in our model are separated by the horizon located at $`u=\beta `$, i.e the singularity of $`\varphi _0`$. By normalizing $`\tau `$ by the normalization factor $`\sqrt{\frac{12}{\mathrm{\Lambda }}}\beta `$ and defining a new coordinate $`r=\mathrm{sinh}\rho `$, one at the end of the day obtains,
$$ds^2=d\tau ^2+a(\tau )^2\left(\frac{dr^2}{1+r^2}+r^2dz_i^2\right),$$
(26)
in which $`a(\tau )=\sqrt{\frac{3}{\mathrm{\Lambda }}}\mathrm{sinh}\sqrt{\frac{\mathrm{\Lambda }}{3}}\tau `$. This is the Robertson-Walker metric for open de Sitter universe. One can easily calculate the energy density $`\rho `$ and the pressure $`p`$ of the cosmological stuff corresponding to $`\varphi _0`$ using the Friedmann equations for the open universe,
$`\left({\displaystyle \frac{\dot{a}}{a}}\right)^2`$ $`=`$ $`{\displaystyle \frac{8\pi G}{3}}\rho +{\displaystyle \frac{1}{a^2}},`$
$`{\displaystyle \frac{\ddot{a}}{a}}`$ $`=`$ $`{\displaystyle \frac{4\pi G}{3}}(\rho +3p).`$ (27)
One verifies that $`p`$ and $`\rho `$ satisfy the equation of state for the cosmological constant $`\rho =p=\mathrm{\Lambda }`$ ($`8\pi G=1`$).
## Acknowledgement
The author gratefully thanks A. Kusenko and V. A. Rubakov for useful discussions. I also thank ICTP for hospitality while this work was completed. The financial support of Isfahan University of Technology (IUT) is acknowledged. |
warning/0506/cond-mat0506213.html | ar5iv | text | # Multipole ordering in 𝑓-electron systems on the basis of a 𝑗-𝑗 coupling scheme
## I Introduction
It is one of currently important issues in the research field of condensed matter physics to unveil exotic magnetic properties of strongly correlated electron materials with active orbital degree of freedom. Among those materials, in $`d`$-electron systems such as transition metal oxides, origin of complex magnetic structure has been vigorously discussed based on the concept of orbital ordering.Imada ; Tokura ; Dagotto ; Hotta0 Also in $`f`$-electron materials including rare-earth and actinide elements, various kinds of magnetic and orbital ordering have been found.Santini2 ; orbital2001 It is now widely recognized that orbital degree of freedom plays a crucial role for the emergence of novel magnetism in $`d`$\- and $`f`$-electron systems.
Here we should note that in $`f`$-electron systems, spin and orbital are not independent degrees of freedom, since they are tightly coupled with each other due to the strong spin-orbit interaction. Then, in order to describe such a complicated spin-orbital coupled system, we usually represent the $`f`$-electron state in terms of “multipole” degree of freedom, rather than using spin and orbital degrees of freedom as in $`d`$-electron systems. Among multipole moments, there have been intensive and extensive studies on dipole and/or quadrupole ordering in $`f`$-electron systems. In usual cases, magnetic ordering indicates dipole one, which can be detected by neutron diffraction experiments. Ordinary orbital ordering means quadrupole one, which can also be detected experimentally, since it induces lattice distortions due to the spatial anisotropy in charge distribution.
In addition to dipole and quadrupole ordering, in recent years, possibility of higher-order multipole ordering, i.e., magnetic octupole ordering, has been also discussed for Ce<sub>x</sub>La<sub>1-x</sub>B<sub>6</sub> Kuramoto ; Kusunose ; Sakakibara ; Kubo ; Kobayashi ; Iwasa ; Suzuki and NpO<sub>2</sub>, Santini ; Paixao ; Lovesey ; Kiss ; Tokunaga ; Sakai2 ; Kubo:NpO2 to reconcile experimental observations which seem to contradict one another at first glance. Very recently, a possibility of octupole ordering has been proposed also for SmRu<sub>4</sub>P<sub>12</sub>Yoshizawa ; Hachitani It is noted that in these materials, crystalline electric field (CEF) ground states are $`\mathrm{\Gamma }_8`$ quartets with large degeneracy even under a CEF potential. Zirngiebl ; Fournier ; Matsuhira In the $`\mathrm{\Gamma }_8`$ ground-state multiplet, octupoles exist as independent moments besides dipole and quadrupole moments.Shiina Then, phenomenological theories have been developed under the assumption that octupole ordering occurs. Note that direct detection of octupole ordering is very difficult, since the octupole moment directly couples to neither a magnetic field nor lattice distortions. However, those phenomenological theories have been successful in explaining several experimental facts consistently, e.g., induced quadrupole moments in octupole ordered states in Ce<sub>x</sub>La<sub>1-x</sub>B<sub>6</sub> Kusunose ; Kubo and NpO<sub>2</sub>.Paixao ; Kiss
As mentioned above, thus far, the study on multipole ordering in $`f`$-electron systems has been almost limited in the phenomenological level, mainly due to the complexity in the treatment of multipole degree of freedom. It might be possible to consider a Heisenberg-like model for multipole moments, but the interactions among multipole moments were determined just phenomenologically. It is highly required to proceed to microscopic theory, in order to understand the origin of multipole ordering in $`f`$-electron systems. However, it is very hard and practically impossible to study multipole ordering in the model retaining all the $`f`$-electron states. Then, it is necessary to consider a tractable model which keeps correct $`f`$-electron symmetry.
One way for such model construction is to use an $`LS`$ coupling scheme. For instance, the Ruderman-Kittel-Kasuya-Yosida (RKKY) interactions were estimated in DyZn, Schmitt and in CeB<sub>6</sub> and CeB<sub>2</sub>C<sub>2</sub> Sakurai from microscopic models using the $`LS`$ coupling scheme. However, the method based on the $`LS`$ coupling scheme is complicated and seems still hard to be extended. One reason for the difficulty is that we cannot apply standard quantum-field theoretical techniques in the $`LS`$ coupling scheme, since Wick’s theorem does not hold. From this viewpoint, it is recommended to use a $`j`$-$`j`$ coupling scheme.Hotta Since individual $`f`$-electron states are first defined, we can include many-body effects in systematic ways using theoretical techniques developed for the research of $`d`$-electron systems. Kubo:NpO2 ; Hotta ; Hotta2 ; Hotta3a ; Hotta3
In this paper, in order to investigate how multipole ordering appears in $`f`$-electron systems from a microscopic viewpoint, we exploit the $`j`$-$`j`$ coupling scheme. We construct tight-binding models on three kinds of lattices, simple cubic (sc), bcc, and fcc, by including Coulomb interactions among $`\mathrm{\Gamma }_8`$ states. In order to discuss multipole ordering in these models, we derive an effective multipole interaction model in the strong-coupling limit for each lattice structure by using the second-order perturbation theory with respect to $`f`$-$`f`$ hopping integrals, as to estimate the superexchange interaction in $`d`$-electron systems. Then, within a mean-field approximation, we clarify what kind of multipole ordering occurs in the effective model: For the sc lattice, a $`\mathrm{\Gamma }_{3g}`$ antiferro-quadrupole transition occurs, while for the bcc lattice, $`\mathrm{\Gamma }_{2u}`$ antiferro-octupole ordering appears. For the fcc lattice with geometrical frustration, we find longitudinal triple-$`𝒒`$ $`\mathrm{\Gamma }_{5u}`$ octupole ordering.
The organization of this paper is as follows. In Sec. II, we introduce a tight-binding model based on the $`j`$-$`j`$ coupling scheme including only the $`\mathrm{\Gamma }_8`$ states. In Sec. III, we describe the general prescription to derive an effective Hamiltonian from the $`\mathrm{\Gamma }_8`$ model. In Sec. IV, we show the mean-field results of the effective models on sc, bcc, and fcc lattices. Finally, in Sec. V, the paper is summarized.
## II Hamiltonian
When we study theoretically the $`f`$-electron properties, the $`LS`$ coupling scheme has been frequently used to include the effect of Coulomb interactions, spin-orbit coupling, and CEF potential. However, as mentioned above, it is not possible to apply standard quantum-field theoretical technique in the $`LS`$ coupling scheme, since Wick’s theorem does not hold. In order to overcome such a difficulty, it has been proposed to construct a microscopic model for $`f`$-electron systems by exploiting the $`j`$-$`j`$ coupling scheme,Hotta where we include first the spin-orbit coupling so as to define the state labeled by the total angular momentum $`j`$. For $`f`$ orbitals with angular momentum $`\mathrm{}`$=3, we immediately obtain an octet with $`j`$=7/2(=3+1/2) and a sextet with $`j`$=5/2(=3$``$1/2), which are well separated by the spin-orbit interaction. Since the spin-orbital coupling is, at least, in the order of 0.1 eV for $`f`$ electrons, it is enough to take into account the $`j`$=5/2 sextet, when we investigate low-temperature properties of $`f`$-electron compounds in the $`j`$-$`j`$ coupling scheme.
In order to construct the many-body state, we accommodate $`f`$ electrons in the $`j`$=5/2 sextet by following the Hund’s rule interactions and CEF potential, as we have done for $`d`$-electron systems. It has been found that the many-electron state obtained in the $`j`$-$`j`$ coupling scheme is continuously changed to the corresponding state in the $`LS`$ coupling scheme, as long as those states in both schemes belong to the same symmetry group.Hotta3 Namely, if we based on the spirit of adiabatic continuation, there is no serious difference between the states of the $`LS`$ and $`j`$-$`j`$ coupling schemes. Depending on the problem, we can use one of the schemes for $`f`$-electron systems. For instance, if we attempt to explain phenomenologically the experimental results of $`f`$-electron insulators, it is highly recommended to use the $`LS`$ coupling scheme. On the other hand, the $`j`$-$`j`$ coupling scheme is rather appropriate to develop a microscopic theory for novel magnetism and unconventional superconductivity of $`f`$-electron systems. In the present paper, our purpose is to construct a microscopic theory for multipole ordering from the viewpoint of spin-orbital complex. Thus, we exploit the $`j`$-$`j`$ coupling scheme throughout this paper.
As described above, we consider only the states with $`j`$=5/2. The $`j`$=5/2 states are further split into $`\mathrm{\Gamma }_7`$ doublet and $`\mathrm{\Gamma }_8`$ quartet due to a cubic CEF. In order to consider multipole phenomena such as octupole ordering in $`f`$-electron systems from a microscopic viewpoint, in this paper we consider only $`\mathrm{\Gamma }_8`$ states by assuming large CEF splitting energy between $`\mathrm{\Gamma }_7`$ and $`\mathrm{\Gamma }_8`$ levels. This simplification is motivated by the fact that the possibility of exotic octupole ordering has been actively discussed in Ce<sub>x</sub>La<sub>1-x</sub>B<sub>6</sub> and NpO<sub>2</sub> with $`\mathrm{\Gamma }_8`$ ground state.
Here readers may be doubtful of the reality of our assumption, since the Coulomb interaction among $`f`$ electrons is naively thought to be larger than the CEF level splitting in any case. However, it should be noted that we are now considering the $`f`$-electron state in the $`j`$-$`j`$ coupling scheme, not in the original $`f`$-electron state with angular momentum $`\mathrm{}`$=3. As pointed out in Ref. Hotta, , the Hund’s rule interaction in the $`j`$-$`j`$ coupling scheme is effectively reduced to be 1/49 of the original Hund’s rule coupling. Namely, even if the original Hund’s rule coupling among $`f`$ electrons is 1 eV, it is reduced to 200 K in the $`j`$-$`j`$ coupling scheme. We note that the CEF level splitting in actinide dioxides is considered to be larger than 1000 K. Kubo:NpO2 ; Kubo:fp We also recall that the CEF level splitting in CeB<sub>6</sub> is as large as 500 K. Zirngiebl Thus, we safely conclude that our present assumption is correctly related to the realistic situation. Of course, in order to achieve quantitative agreement with experimental results, it is necessary to include also $`\mathrm{\Gamma }_7`$ level, since the magnitude of the CEF splitting is always finite, even if it is large compared with the effective Hund’s rule interaction. However, we strongly believe that it is possible to grasp microscopic origin of multipole ordering in $`f`$-electron systems on the basis of the $`\mathrm{\Gamma }_8`$ model, since this model is considered to be connected adiabatically from the realistic situation. We postpone further effort to develop more general theory to include all the $`j`$=5/2 sextet in future.
Concerning the $`f`$-electron number, in this paper we treat only the case with one $`f`$ electron in the $`\mathrm{\Gamma }_8`$ multiplet per site. However, this restriction does $`not`$ simply indicate that we consider only the Ce-based compound. In the $`j`$-$`j`$ coupling scheme, in order to consider $`f^n`$-electron systems, where $`n`$ indicates local $`f`$ electron number per site, we accommodate $`f`$ electrons in the one-electron CEF levels due to the balance between Coulomb interactions and CEF level splitting energy, just as in the case of $`d`$-electron systems. Thus, the situation with one $`f`$ electron in the $`\mathrm{\Gamma }_8`$ multiplet per site expresses both cases with $`n`$=1 in the $`\mathrm{\Gamma }_8`$-$`\mathrm{\Gamma }_7`$ \[Fig. 1(a)\] and $`n`$=3 in the $`\mathrm{\Gamma }_7`$-$`\mathrm{\Gamma }_8`$ \[Fig. 1(b)\] systems, where $`\mathrm{\Gamma }_x`$-$`\mathrm{\Gamma }_y`$ symbolically denotes the situation with $`\mathrm{\Gamma }_x`$ ground and $`\mathrm{\Gamma }_y`$ excited states. Furthermore, we should note that due to the electron-hole symmetry in the $`\mathrm{\Gamma }_8`$ subspace, the effective model with one $`f`$ electron in the $`\mathrm{\Gamma }_8`$ state is the same for that in the case with three electrons in the $`\mathrm{\Gamma }_8`$ multiplet. Namely, the present model also indicates both cases with $`n`$=3 in the $`\mathrm{\Gamma }_8`$-$`\mathrm{\Gamma }_7`$ \[Fig. 1(c)\] and $`n`$=5 in the $`\mathrm{\Gamma }_7`$-$`\mathrm{\Gamma }_8`$ \[Fig. 1(d)\] systems.
Before proceeding to the exhibition of the Hamiltonian, it is necessary to define $`f`$-electron operators in $`\mathrm{\Gamma }_8`$ states. Since the $`\mathrm{\Gamma }_8`$ quartet consists of two Kramers doublets, we introduce orbital index $`\tau `$ (=$`\alpha `$ and $`\beta `$) to distinguish the two Kramers doublets, while spin index $`\sigma `$ (=$``$ and $``$) is defined to distinguish the two states in each Kramers doublet. In the second-quantized form, annihilation operators for $`\mathrm{\Gamma }_8`$ electrons are defined as
$`f_{𝐫\alpha }`$ $`=\sqrt{5/6}a_{𝐫5/2}+\sqrt{1/6}a_{𝐫3/2},`$ (1a)
$`f_{𝐫\alpha }`$ $`=\sqrt{5/6}a_{𝐫5/2}+\sqrt{1/6}a_{𝐫3/2},`$ (1b)
for $`\alpha `$-orbital electrons, and
$`f_{𝐫\beta }`$ $`=a_{𝐫1/2},`$ (2a)
$`f_{𝐫\beta }`$ $`=a_{𝐫1/2},`$ (2b)
for $`\beta `$-orbital electrons, where $`a_{𝐫j_z}`$ is the annihilation operator for an electron with the $`z`$-component $`j_z`$ of the total angular momentum at site $`𝐫`$.
Now we show the Hamiltonian of $`\mathrm{\Gamma }_8`$ electrons. For the purpose to consider the effective model later, it is convenient to express the Hamiltonian in the form of
$$=_{\text{kin}}+_{\text{loc}},$$
(3)
where $`_{\mathrm{kin}}`$ denotes the kinetic term of $`f`$ electrons and $`_{\text{loc}}`$ indicates the local interaction part for $`\mathrm{\Gamma }_8`$ electrons. In this paper, the kinetic term of $`\mathrm{\Gamma }_8`$ electrons is given by exploiting the tight-binding approximation. Then, $`_{\text{kin}}`$ is expressed as
$$_{\text{kin}}=\underset{𝐫,𝝁,\tau ,\sigma ,\tau ^{},\sigma ^{}}{}t_{\tau \sigma ;\tau ^{}\sigma ^{}}^𝝁f_{𝐫\tau \sigma }^{}f_{𝐫+𝝁\tau ^{}\sigma ^{}},$$
(4)
where $`𝝁`$ is a vector connecting nearest-neighbor sites and $`t_{\tau \sigma ;\tau ^{}\sigma ^{}}^𝝁`$ is the hopping integral of an electron with $`(\tau ^{},\sigma ^{})`$ at site $`𝐫`$+$`𝝁`$ to the $`(\tau ,\sigma )`$ state at $`𝐫`$. We note that the hopping integral $`t_{\tau \sigma ;\tau ^{}\sigma ^{}}^𝝁`$ depends on orbital, spin, and direction $`𝝁`$, due to $`f`$-electron symmetry. Then, the form of the hopping integral is characteristic of lattice structure. The explicit form of the hopping matrix will be shown later for each lattice structure. Note also the relation $`t_{\tau \sigma ;\tau ^{}\sigma ^{}}^𝝁`$ =$`t_{\tau \sigma ;\tau ^{}\sigma ^{}}^𝝁`$.
As for the local $`f`$-electron term $`_{\mathrm{loc}}`$, since we assume the large CEF splitting energy between $`\mathrm{\Gamma }_7`$ and $`\mathrm{\Gamma }_8`$ levels, it is enough to consider the Coulomb interaction terms among $`\mathrm{\Gamma }_8`$ electrons. As easily understood from the introduction of ‘spin’ and ‘orbital’ in the $`j`$-$`j`$ coupling scheme, the local $`f`$-electron term in the $`\mathrm{\Gamma }_8`$ quartet becomes the same as that of the two-orbital systems for $`d`$ electrons. In fact, after lengthy algebraic calculations for Racah parameters in the $`j`$-$`j`$ coupling scheme,Hotta $`_{\mathrm{loc}}`$ is given as
$$\begin{array}{cc}\hfill _{\text{loc}}& =U\underset{𝐫\tau }{}n_{𝐫\tau }n_{𝐫\tau }+U^{}\underset{𝐫}{}n_{𝐫\alpha }n_{𝐫\beta }\hfill \\ & +J\underset{𝐫,\sigma ,\sigma ^{}}{}f_{𝐫\alpha \sigma }^{}f_{𝐫\beta \sigma ^{}}^{}f_{𝐫\alpha \sigma ^{}}f_{𝐫\beta \sigma }\hfill \\ & +J^{}\underset{𝐫,\tau \tau ^{}}{}f_{𝐫\tau }^{}f_{𝐫\tau }^{}f_{𝐫\tau ^{}}f_{𝐫\tau ^{}},\hfill \end{array}$$
(5)
where $`n_{𝐫\tau \sigma }`$ =$`f_{𝐫\tau \sigma }^{}f_{𝐫\tau \sigma }`$ and $`n_{𝐫\tau }`$=$`_\sigma n_{𝐫\tau \sigma }`$. The coupling constants $`U`$, $`U^{}`$, $`J`$, and $`J^{}`$ denote the intra-orbital Coulomb, inter-orbital Coulomb, exchange, and pair-hopping interactions, respectively. These are expressed in terms of Racah parameters, and we obtain the relation $`U`$=$`U^{}`$+$`J`$+$`J^{}`$, which can be understood from the rotational invariance in orbital space.Hotta Note that for $`d`$-electron systems, one also has the relation $`J`$=$`J^{}`$. When the electronic wave-function is real, this relation is easily demonstrated from the definition of the Coulomb integral. However, in the $`j`$-$`j`$ coupling scheme the wave-function is complex, and $`J`$ is not equal to $`J^{}`$ in general.
## III Effective Model
In this section, we describe a method to derive an effective Hamiltonian by using the second-order perturbation theory with respect to hopping integrals. Here we emphasize that the procedure is essentially the same as to estimate superexchange interactions for $`d`$-electron systems, although the calculations are tedious due to the existence of orbital degree of freedom. After that, we will apply the standard mean-field theory to the effective model to depict the phase diagram including multipole ordered states. We believe that it is meaningful to understand the complicated $`f`$-electron multipole problem by using a simple $`d`$-electron-like procedure and approximations, both from conceptual and practical viewpoints.
When the Hamiltonian is written in the form of Eq. (3), first we solve the local problem
$$_{\mathrm{loc}}|\mathrm{\Phi }_n^a=E_n|\mathrm{\Phi }_n^a,$$
(6)
where $`E_n`$ denotes the $`n`$-th eigenenergy and $`|\mathrm{\Phi }_n^a`$ is the corresponding eigenstate with a label $`a`$ to distinguish the degenerate states. Since we accommodate one electron per site, the ground state $`|\mathrm{\Phi }_0^a`$ is expressed as
$$|\mathrm{\Phi }_0^a=\underset{𝐫,\tau ,\sigma }{}f_{𝐫\tau \sigma }^{P_a(𝐫,\tau ,\sigma )}|0,$$
(7)
where $`|0`$ is the vacuum state of $`f`$ electrons, $`a`$ denotes the electron configuration in the ground state with one electron per site, and $`P_a(𝐫,\tau ,\sigma )`$ takes 0 or 1 depending on the configuration $`a`$.
Here we consider the formal perturbation expansion in terms of $`_{\mathrm{kin}}`$ in order to construct the effective model. Within the second order, $`_{\mathrm{eff}}`$ is generally written as
$$_{\mathrm{eff}}=\underset{a,b,u}{}\underset{m0}{}|\mathrm{\Phi }_0^a\mathrm{\Phi }_0^a|_{\mathrm{kin}}\frac{|\mathrm{\Phi }_m^u\mathrm{\Phi }_m^u|}{E_0E_m}_{\mathrm{kin}}|\mathrm{\Phi }_0^b\mathrm{\Phi }_0^b|,$$
(8)
where $`a`$ and $`b`$ are labels to distinguish the ground states, while $`u`$ is the label for the degenerate excited states.
Since we consider the situation with one electron per site, the intermediate state due to one-electron hopping has a vacant and a double occupied site. The double occupied site has six possible states, composed of $`\mathrm{\Gamma }_5`$ triplet with energy $`U^{}`$$``$$`J`$, $`\mathrm{\Gamma }_3`$ doublet with energy $`U^{}`$+$`J`$(=$`U`$$``$$`J^{}`$), and $`\mathrm{\Gamma }_1`$ singlet with energy $`U`$+$`J^{}`$. For the mathematical completion, it is necessary to include all possible excited states in the intermediate process, but the calculation becomes complicated. Then, in this paper, in order to grasp the essential point of the $`\mathrm{\Gamma }_8`$ model by avoiding tedious calculations, we include only the lowest-energy $`\mathrm{\Gamma }_5`$ triplet among the intermediate $`f^2`$ states. This restriction to the intermediate states is validated, when $`J`$ is much larger than the hopping energy of $`f`$ electron. Since the $`f`$-electron hopping amplitude is considered to be small compared with $`J`$, even if we also include the hybridization with conduction electrons, this approximation is acceptable in $`f`$-electron systems.
Let us explain the prescription to derive the effective model in the present case. It is convenient to consider exchange processes of electrons between two sites, $`𝐫`$ and $`𝐫^{}`$. Since we consider the situation with one $`f`$ electron per site, the initial state $`|𝐫s_1,𝐫^{}s_2`$ is written as
$$|𝐫s_1,𝐫^{}s_2=f_{𝐫s_1}^{}f_{𝐫^{}s_2}^{}|0,$$
(9)
where $`s_1`$ and $`s_2`$ symbolically denote spin and orbital states for both electrons. Then, we move one electron from the site $`𝐫^{}`$ to $`𝐫`$. As mentioned above, the intermediate $`f^2`$ states at the site $`𝐫`$ is restricted only as the lowest-energy $`\mathrm{\Gamma }_5`$ triplet states. Namely, the intermediated states should be expressed as $`|u,𝐫`$ with the label $`u`$ to distinguish the triplet states, given by
$`|+1,𝐫`$ $`=f_{𝐫\alpha }^{}f_{𝐫\beta }^{}|0,`$ (10a)
$`|0,𝐫`$ $`=(f_{𝐫\alpha }^{}f_{𝐫\beta }^{}+f_{𝐫\alpha }^{}f_{𝐫\beta }^{})|0/\sqrt{2},`$ (10b)
$`|1,𝐫`$ $`=f_{𝐫\alpha }^{}f_{𝐫\beta }^{}|0.`$ (10c)
In order to obtain the effective model Eq. (8), it is enough to evaluate the inner product
$$P_{u;s,s^{}}=𝐫,u|𝐫s,𝐫s^{}.$$
(11)
This quantity is explicitly given by $`P_{+1;\alpha \beta }`$=1, $`P_{0;\alpha \beta }`$=$`1/\sqrt{2}`$, and the other non-zero elements are given by $`P_{u;s^{}s}`$=$`P_{u;ss^{}}`$ and $`P_{u;\tau \sigma \tau ^{}\sigma ^{}}`$ =$`P_{u;\tau \sigma \tau ^{}\sigma ^{}}`$.
Then, by including the processes in which an electron at $`𝐫`$ moves first, we obtain the effective Hamiltonian as
$$_{\mathrm{eff}}=\underset{𝐫,𝐫^{}}{}\underset{s_1\text{}s_4}{}I_{s_3,s_4;s_1,s_2}^{𝐫^{}𝐫}f_{𝐫s_3}^{}f_{𝐫s_1}f_{𝐫^{}s_4}^{}f_{𝐫^{}s_2},$$
(12)
where $`𝐫,𝐫^{}`$ denotes the pair of nearest-neighbor sites and the generalized exchange interaction $`I`$ is given by
$$\begin{array}{cc}\hfill I_{s_3,s_4;s_1,s_2}^{𝐫^{}𝐫}& =\underset{u,s,s^{}}{}[(t_{s^{};s_4}^{𝐫^{}𝐫})^{}P_{u;s_3,s^{}}^{}P_{u;s_1,s}t_{s;s_2}^{𝐫^{}𝐫}\hfill \\ & +(t_{s^{};s_3}^{𝐫𝐫^{}})^{}P_{u;s_4,s^{}}^{}P_{u;s_2,s}t_{s;s_1}^{𝐫𝐫^{}}]/(U^{}J).\hfill \end{array}$$
(13)
In order to investigate the multipole ordering, it is more convenient to express the effective Hamiltonian Eq. (12) in terms of multipole operators. For the purpose, we introduce some notations to describe multipole operators as
$`\stackrel{~}{1}_{\tau \sigma ;\tau ^{}\sigma ^{}}`$ $`\delta _{\tau \tau ^{}}\delta _{\sigma \sigma ^{}},`$ (14a)
$`\stackrel{~}{𝝉}_{\tau \sigma ;\tau ^{}\sigma ^{}}`$ $`𝝈_{\tau \tau ^{}}\delta _{\sigma \sigma ^{}},`$ (14b)
$`\stackrel{~}{𝝈}_{\tau \sigma ;\tau ^{}\sigma ^{}}`$ $`\delta _{\tau \tau ^{}}𝝈_{\sigma \sigma ^{}},`$ (14c)
$`\stackrel{~}{\eta }^\pm `$ $`(\pm \sqrt{3}\stackrel{~}{\tau }^x\stackrel{~}{\tau }^z)/2,`$ (14d)
$`\stackrel{~}{\xi }^\pm `$ $`(\stackrel{~}{\tau }^x\pm \sqrt{3}\stackrel{~}{\tau }^z)/2,`$ (14e)
where $`𝝈`$ are the Pauli matrices. By using these notations, we define one-particle operators at site $`𝐫`$ as
$$\widehat{A}_𝐫\underset{\tau \tau ^{}\sigma \sigma ^{}}{}f_{𝐫\tau \sigma }^{}\stackrel{~}{A}_{\tau \sigma ;\tau ^{}\sigma ^{}}f_{𝐫\tau ^{}\sigma ^{}},$$
(15)
where $`\stackrel{~}{A}`$ is a $`4\times 4`$ matrix. The multipole operators in the $`\mathrm{\Gamma }_8`$ subspace are listed in Table 1.
With the use of above multipole operators, the effective Hamiltonian is finally arranged in the form of
$$_{\text{eff}}=\underset{𝐪}{}(_{1𝐪}+_{2𝐪}+_{4u1𝐪}+_{4u2𝐪}),$$
(16)
where $`𝐪`$ is the wave vector and $`_{1𝐪}`$ denotes quadrupole interactions. $`_{4un𝐪}`$ ($`n`$=1 or 2) denotes interactions between $`\mathrm{\Gamma }_{4un}`$ moments and ones between $`\mathrm{\Gamma }_{4un}`$ and other octupole moments with symmetry different from $`\mathrm{\Gamma }_{4u}`$. $`_{2𝐪}`$ denotes other dipole and octupole interactions. In general, $`_{2𝐪}`$ includes interactions between $`\mathrm{\Gamma }_{4u1}`$ and $`\mathrm{\Gamma }_{4u2}`$ moments, but we find that such interactions are not included in the models with hopping integrals only through $`(ff\sigma )`$ bonding on sc, bcc, and fcc lattices. The explicit form of each multipole interaction sensitively depends on the lattice structure, as shown in the next section.
## IV Results
Now we can calculate the effective interaction between two electrons located along any direction $`𝐫^{}𝐫`$ by using Eq. (13), if the hopping integral along this direction is determined. The hopping integrals of $`f`$ electrons are evaluated by using the Slater-Koster table. Takegahara In this section, we consider the nearest-neighbor hopping integrals through $`(ff\sigma )`$ bonding for three lattice structures, sc, bcc, and fcc. Then, we present the effective Hamiltonian and its ordered states for each lattice. The structure of our effective model is consistent with the general form of nearest-neighbor multipole interactions on each lattice derived by Sakai et al. Sakai We follow the notation in Ref. Sakai, for convenience.
### IV.1 sc lattice
The nearest-neighbor hopping integrals through $`(ff\sigma )`$ bonding for the sc lattice are given by
$`t^{(a,0,0)}`$ $`=[\stackrel{~}{1}\stackrel{~}{\eta }^+]t_1,`$ (17a)
$`t^{(0,a,0)}`$ $`=[\stackrel{~}{1}\stackrel{~}{\eta }^{}]t_1,`$ (17b)
$`t^{(0,0,a)}`$ $`=[\stackrel{~}{1}\stackrel{~}{\tau }^z]t_1,`$ (17c)
where $`a`$ is the lattice constant and $`t_1`$=$`3(ff\sigma )/14`$.
For the sc lattice, the quadrupole interaction term in Eq. (16) is given by
$$_{1𝐪}=a_1(O_{2,𝐪}^0O_{2,𝐪}^0C_z+\text{c.p.}),$$
(18)
where c.p. denotes cyclic permutations and $`C_\nu `$=$`\mathrm{cos}(q_\nu a)`$ ($`\nu `$=$`x`$, $`y`$, or $`z`$). The value of the coupling constant $`a_1`$ is given in Table 2.
Note that $`O_{2𝐪}^0`$ transforms to $`(\sqrt{3}O_{2𝐪}^2O_{2𝐪}^0)/2`$ and $`(\sqrt{3}O_{2𝐪}^2O_{2𝐪}^0)/2`$ under c.p. $`(x,y,z)(y,z,x)`$ and $`(x,y,z)(z,x,y)`$, respectively. The dipole and octupole interactions are given by
$$_{2𝐪}=b_6[T_{z,𝐪}^{5u}T_{z,𝐪}^{5u}(C_x+C_y)+\text{c.p.}],$$
(19)
and
$$\begin{array}{cc}\hfill _{4un𝐪}=b_1^{(n)}& [J_{z𝐪}^{4un}J_{z𝐪}^{4un}C_z+\text{c.p.}]\hfill \\ \hfill +b_2^{(n)}& [J_{z𝐪}^{4un}J_{z𝐪}^{4un}(C_x+C_y)+\text{c.p.}]\hfill \\ \hfill +b_3^{(n)}& [T_{z𝐪}^{5u}J_{z𝐪}^{4un}(C_xC_y)+\text{c.p.})],\hfill \end{array}$$
(20)
where values of the coupling constants $`b_i`$ and $`b_i^{(n)}`$ are shown in Table 2.
Note that the form of the hopping integrals Eqs. (17) are exactly the same as those for the $`e_g`$ orbitals of $`d`$ electrons via $`(dd\sigma )`$ bonding. Hotta ; Anderson Thus, the effective Hamiltonian has the same form as in the $`e_g`$ model considering only the lowest-energy intermediate states, Kugel when we interpret that $`\tau `$ and $`\sigma `$ denote $`e_g`$ orbital and real spin, respectively. However, the physical meaning of the present model is different from that of the $`e_g`$ model. In particular, the effect of a magnetic field is essentially different. The dipole moment which couples to a magnetic field $`𝐇`$ is given by $`𝐉`$=$`(7/6)[𝐉^{4u1}`$+$`(4/7)𝐉^{4u2}]`$ for the $`\mathrm{\Gamma }_8`$ model, while for the $`e_g`$ model, real spin $`𝝈`$ of $`d`$ electrons is simply coupled to a magnetic field. In contrast to the $`e_g`$ model, a magnetic field resolves the degeneracy in the $`\tau `$ space even within a mean-field theory for the present model, as we will see later.
By applying mean-field theory to the effective model, we find a $`\mathrm{\Gamma }_{3g}`$ antiferro-quadrupole transition at a temperature $`T`$=$`T_{3g}`$=$`3a_1/k_\text{B}`$. As lowering temperature further, we find a $`\mathrm{\Gamma }_{4u1}`$ ferromagnetic transition. This ferromagnetic transition can be regarded as a $`\mathrm{\Gamma }_{5u}`$ antiferro-octupole transition, since the $`\mathrm{\Gamma }_{4u1}`$ ferromagnetic state with the $`\mathrm{\Gamma }_{3g}`$ antiferro-quadrupole moment is equivalent to the $`\mathrm{\Gamma }_{5u}`$ antiferro-octupole ordered state with the $`\mathrm{\Gamma }_{3g}`$ antiferro-quadrupole moment. The ground state energy is $`(3/2)a_12b_6+b_1^{(1)}+2b_2^{(1)}`$ per site.
In Fig. 2(a), we depict an $`H`$-$`T`$ phase diagram. We note that the ferromagnetic transition at zero magnetic field turns to be a crossover under the finite magnetic field. The crossover is drawn by dashed curve, determined by the peak position in the magnetic susceptibility. Since it is found that the crossover curve is almost isotropic in the region shown here, we depict only the curve for $`𝐇[001]`$. Note also that under a magnetic field, $`\mathrm{\Gamma }_{4u1}`$ moments become finite, and then, the $`\mathrm{\Gamma }_{5u}`$ antiferro-octupole interaction ($`b_6`$$`>`$0) effectively becomes a $`\mathrm{\Gamma }_{3g}`$ antiferro-quadrupole interaction. Thus, the $`\mathrm{\Gamma }_{3g}`$ antiferro-quadrupole transition temperature increases as $`H`$ is increased at a low magnetic field region. This behavior reminds us of the experimental results for CeB<sub>6</sub>, although the order parameter in the quadrupole ordered phase of CeB<sub>6</sub> is the $`\mathrm{\Gamma }_{5g}`$ quadrupole moment. Magnetization as a function of $`H`$ is shown in Fig. 2(b). The magnetization is isotropic as $`H`$$``$0 since the $`\mathrm{\Gamma }_{4u1}`$ moment is isotropic, while anisotropy develops under a high magnetic field.
In Figs. 3(a)–(c), we show specific heat, magnetization, and magnetic susceptibility, respectively, as functions of temperature. We observe two-step jump of specific heat at the quadrupole and ferromagnetic transition temperatures, since we have applied the mean-field theory to these second-order transitions. Note that the magnetization starts to develop below the ferromagnetic transition temperature. The magnetic susceptibility exhibits a bend at $`T_{3g}`$, while it diverges at the ferromagnetic transition temperature. Under the magnetic field, this divergence turns to be a peak, which defines the crossover to the ferromagnetic state in the $`H`$-$`T`$ phase diagram.
Without magnetic field, the orbital ($`\tau `$) state is continuously degenerate in the mean-field theory, although such continuous symmetry is absent in this model. As has been discussed for an $`e_g`$ electron model such as perovskite manganites,Brink quantum fluctuations can resolve this continuous degeneracy, but in the present model with the strong spin-orbit interaction, magnetic field can resolve this degeneracy. The ground states are ferromagnetic with $`𝐉^{4u1}𝐇`$, where $`\mathrm{}`$ denotes the expectation value. Accompanied $`O_2^2`$ ordering is G-type \[$`𝐪=(1/2,1/2,1/2)`$ in units of $`2\pi /a`$\] or C-type \[$`𝐪=(1/2,1/2,0)`$\] for $`𝐇[001]`$, while for $`𝐇[110]`$, it is C-type. For $`𝐇[111]`$, there appear C-type $`O_2^2`$ ordering or equivalent ones in the cubic symmetry.
### IV.2 bcc lattice
The hopping integrals for the bcc lattice are given by
$`t^{(a/2,a/2,a/2)}`$ $`=[\stackrel{~}{1}+\stackrel{~}{\tau }^y(+\stackrel{~}{\sigma }^x+\stackrel{~}{\sigma }^y+\stackrel{~}{\sigma }^z)/\sqrt{3}]t_2,`$ (21a)
$`t^{(a/2,a/2,a/2)}`$ $`=[\stackrel{~}{1}+\stackrel{~}{\tau }^y(+\stackrel{~}{\sigma }^x\stackrel{~}{\sigma }^y\stackrel{~}{\sigma }^z)/\sqrt{3}]t_2,`$ (21b)
$`t^{(a/2,a/2,a/2)}`$ $`=[\stackrel{~}{1}+\stackrel{~}{\tau }^y(\stackrel{~}{\sigma }^x+\stackrel{~}{\sigma }^y\stackrel{~}{\sigma }^z)/\sqrt{3}]t_2,`$ (21c)
$`t^{(a/2,a/2,a/2)}`$ $`=[\stackrel{~}{1}+\stackrel{~}{\tau }^y(\stackrel{~}{\sigma }^x\stackrel{~}{\sigma }^y+\stackrel{~}{\sigma }^z)/\sqrt{3}]t_2,`$ (21d)
where $`a`$ is the lattice constant and $`t_2`$=$`2(ff\sigma )/21`$.
After some algebraic calculations, we obtain the quadrupole interaction term for the bcc lattice as
$$\begin{array}{cc}\hfill _{1𝐪}=a_3& (O_{xy,𝐪}O_{xy,𝐪}+\text{c.p.})c_xc_yc_z\hfill \\ \hfill +a_4& [O_{yz,𝐪}O_{zx,𝐪}s_xs_yc_z+\text{c.p.}].\hfill \end{array}$$
(22)
The dipole and octupole interactions are given by
$$\begin{array}{cc}\hfill _{2𝐪}=b_5& T_{xyz,𝐪}T_{xyz,𝐪}c_xc_yc_z\hfill \\ \hfill +b_6& (T_{z,𝐪}^{5u}T_{z,𝐪}^{5u}+\text{c.p.})c_xc_yc_z\hfill \\ \hfill +b_7& [T_{x,𝐪}^{5u}T_{y,𝐪}^{5u}s_xs_yc_z+\text{c.p.}],\hfill \end{array}$$
(23)
and
$$\begin{array}{cc}\hfill _{4un𝐪}=b_1^{(n)}& (J_{z𝐪}^{4un}J_{z𝐪}^{4un}+\text{c.p.})c_xc_yc_z\hfill \\ \hfill +b_2^{(n)}& [J_{x𝐪}^{4un}J_{y𝐪}^{4un}s_xs_yc_z+\text{c.p.}]\hfill \\ \hfill +b_3^{(n)}& T_{xyz𝐪}(J_{z𝐪}^{4un}s_xs_yc_z+\text{c.p.})\hfill \\ \hfill +b_4^{(n)}& [T_{z𝐪}^{5u}(J_{x𝐪}^{4un}s_zs_xc_y+J_{y𝐪}^{4un}s_ys_zc_x)+\text{c.p.}],\hfill \end{array}$$
(24)
where $`c_\nu `$=$`\mathrm{cos}(q_\nu a/2)`$ and $`s_\nu `$=$`\mathrm{sin}(q_\nu a/2)`$. The values of the coupling constants $`a_i`$, $`b_i`$, and $`b_i^{(n)}`$ are shown in Table 3.
In the mean-field approximation, we find a $`\mathrm{\Gamma }_{2u}`$ antiferro-octupole transition at $`T_{2u}`$=$`2b_5/k_\text{B}`$ with $`𝐪`$=$`(1,0,0)`$, and a $`\mathrm{\Gamma }_{4u1}`$ ferromagnetic transition at a lower temperature. The ground state has the $`\mathrm{\Gamma }_{5g}`$ antiferro-quadrupole moment with the same ordering wave-vector as the $`\mathrm{\Gamma }_{2u}`$ moment. The ground state energy is $`a_3b_5+b_1^{(1)}`$ per site.
In Fig. 4(a), we show an $`H`$-$`T`$ phase diagram. Again the ferromagnetic transition becomes a crossover under the finite magnetic field. The crossover curve determined by the peak in the magnetic susceptibility is found to be almost isotropic in the region shown here. Then, we show only the curve for $`𝐇[001]`$. In the region for high $`H`$ and low $`T`$, we find two uniform phases. One is a phase with uniform $`T_{xyz𝐫}`$ depending on temperature and another is a phase with uniform $`T_{xyz𝐫}`$ which does not depend on temperature, as shown in Figs. 4(a) and (c).
In Fig. 4(b), we show magnetization as a function of $`H`$. We note that the magnetization is isotropic as $`H`$$``$0 as in the sc lattice, since the order parameter of the ferromagnetic transition is the $`\mathrm{\Gamma }_{4u1}`$ moment. Note also that the jump in the magnetization at $`(g_J\mu _\text{B}H)/(k_\text{B}T_{2u})`$=5.4 for $`𝐇[111]`$ indicates the transition to the uniform state.
Figures 5(a)–(c) show specific heat, magnetization and magnetic susceptibility as functions of temperature, respectively. We observe two jumps in the specific heat at the octupole and ferromagnetic transition temperatures. The magnetization begins to develop below the ferromagnetic transition temperature. The magnetic susceptibility has a bend at $`T_{2u}`$ and diverges at the ferromagnetic transition temperature. Note that the anomaly in the magnetic susceptibility at $`T_{2u}`$ is very weak. In the pure magnetic $`\mathrm{\Gamma }_{2u}`$ octupole ordered state, there remains degeneracy, while in ordinary magnetic states, degeneracy is fully resolved. Thus, the nature of the $`\mathrm{\Gamma }_{2u}`$ octupole phase is similar to that of the quadrupole ordered phases. For instance, the anomaly in the magnetic susceptibility is weak at the transition temperature, there is no ordered magnetic dipole moment, and another phase transition occurs at a lower temperature.
The ground state is continuously degenerate, since the $`\mathrm{\Gamma }_{4u1}`$ and $`\mathrm{\Gamma }_{5g}`$ moments are isotropic in this model within the $`𝐪=(1,0,0)`$ structure. We note that this degeneracy is due to the symmetry of the model in contrast to the sc lattice. By applying a magnetic field, the ground states are uniquely determined. The ground states are ferromagnetic phases $`𝐉^{4u1}𝐇`$ with antiferro $`O_{xy}`$ ordering for $`𝐇[001]`$, with antiferro $`O_{yz}+O_{zx}`$ ordering for $`𝐇[110]`$, and with antiferro $`O_{yz}+O_{zx}+O_{xy}`$ ordering for $`𝐇[111]`$.
As mentioned in Sec. I, quite recently, a possibility of octupole ordering in filled skutterudite compound SmRu<sub>4</sub>P<sub>12</sub> has been suggested experimentally.Yoshizawa ; Hachitani In the filled skutterudite structure, rare-earth ion surrounded by pnictogens form the bcc lattice. Moreover, the $`\mathrm{\Gamma }_8`$ CEF ground state has been reported in the Sm-based filled skutterudite. Matsuhira Thus, we expect to apply the present model to Sm-based filled skutterudites. When we compare our result on the bcc lattice with the experimental suggestion, octupole ordering actually occurs in our model for the bcc lattice, but $`\mathrm{\Gamma }_{2u}`$ octupole ordered state does not seem to explain the experimental results. This discrepancy is due to the suppression of $`\mathrm{\Gamma }_7`$ orbital, since in the filled skutterudites, conduction electron has $`a_u`$ symmetry, which hybridizes with $`\mathrm{\Gamma }_7`$ electron. In addition, the level splitting between $`\mathrm{\Gamma }_7`$ and $`\mathrm{\Gamma }_8`$ is considered to be rather small in filled skutterudites. Thus, for filled-skutterudite materials, we should consider the $`j`$=5/2 sextet model in the bcc lattice with itinerant $`\mathrm{\Gamma }_7`$ and localized $`\mathrm{\Gamma }_8`$ orbitals. We postpone the analysis of such a model in future.
### IV.3 fcc lattice
The hopping integrals for the fcc lattice are given by
$`t^{(0,a/2,a/2)}`$ $`=[\stackrel{~}{1}+(\stackrel{~}{\eta }^+4\sqrt{3}\stackrel{~}{\tau }^y\stackrel{~}{\sigma }^x)/7]t_3,`$ (25a)
$`t^{(a/2,0,a/2)}`$ $`=[\stackrel{~}{1}+(\stackrel{~}{\eta }^{}4\sqrt{3}\stackrel{~}{\tau }^y\stackrel{~}{\sigma }^y)/7]t_3,`$ (25b)
$`t^{(a/2,a/2,0)}`$ $`=[\stackrel{~}{1}+(\stackrel{~}{\tau }^z4\sqrt{3}\stackrel{~}{\tau }^y\stackrel{~}{\sigma }^z)/7]t_3,`$ (25c)
$`t^{(0,a/2,a/2)}`$ $`=[\stackrel{~}{1}+(\stackrel{~}{\eta }^++4\sqrt{3}\stackrel{~}{\tau }^y\stackrel{~}{\sigma }^x)/7]t_3,`$ (25d)
$`t^{(a/2,0,a/2)}`$ $`=[\stackrel{~}{1}+(\stackrel{~}{\eta }^{}+4\sqrt{3}\stackrel{~}{\tau }^y\stackrel{~}{\sigma }^y)/7]t_3,`$ (25e)
$`t^{(a/2,a/2,0)}`$ $`=[\stackrel{~}{1}+(\stackrel{~}{\tau }^z+4\sqrt{3}\stackrel{~}{\tau }^y\stackrel{~}{\sigma }^z)/7]t_3,`$ (25f)
where $`a`$ is the lattice constant and $`t_3`$=$`(ff\sigma )/8`$.
Each multipole interaction term in the effective Hamiltonian for the fcc lattice is given by
$$\begin{array}{cc}\hfill _{1𝐪}=a_1& (O_{2,𝐪}^0O_{2,𝐪}^0c_xc_y+\text{c.p.})\hfill \\ \hfill +a_3& (O_{2,𝐪}^0O_{xy,𝐪}s_xs_y+\text{c.p.})\hfill \\ \hfill +a_4& (O_{xy,𝐪}O_{xy,𝐪}c_xc_y+\text{c.p.}),\hfill \end{array}$$
(26)
$$\begin{array}{cc}\hfill _{2𝐪}& =b_8[T_{z,𝐪}^{5u}T_{z,𝐪}^{5u}(c_yc_z+c_zc_x)+\text{c.p.}]\hfill \\ & +b_9[T_{x,𝐪}^{5u}T_{y,𝐪}^{5u}s_xs_y+\text{c.p.}]\hfill \\ & +b_{10}T_{xyz,𝐪}T_{xyz,𝐪}(c_xc_y+\text{c.p.}),\hfill \end{array}$$
(27)
and
$$\begin{array}{cc}\hfill _{4un𝐪}=b_1^{(n)}& [J_{z𝐪}^{4un}J_{z𝐪}^{4un}c_xc_y+\text{c.p.}]\hfill \\ \hfill +b_2^{(n)}& [J_{z𝐪}^{4un}J_{z𝐪}^{4un}(c_yc_z+c_zc_x)+\text{c.p.}]\hfill \\ \hfill +b_3^{(n)}& [J_{x𝐪}^{4un}J_{y𝐪}^{4un}s_xs_y+\text{c.p.}]\hfill \\ \hfill +b_4^{(n)}& [T_{xyz𝐪}(J_{z𝐪}^{4un}s_xs_y+\text{c.p.})]\hfill \\ \hfill +b_5^{(n)}& [T_{z𝐪}^{5u}J_{z𝐪}^{4un}c_z(c_xc_y)+\text{c.p.})]\hfill \\ \hfill +b_6^{(n)}& [T_{z𝐪}^{5u}(J_{x𝐪}^{4un}s_zs_x+J_{y𝐪}^{4un}s_ys_z)+\text{c.p.}].\hfill \end{array}$$
(28)
The values of the coupling constants $`a_i`$, $`b_i`$ and $`b_i^{(n)}`$ are shown in Table 4.
As already mentioned in Ref. Kubo:NpO2, , it is necessary to analyze the effective model carefully for the fcc lattice, since the model includes geometrical frustration. It is risky to apply directly the mean-field approximation to the effective model. First we evaluate the correlation function in the ground state using an unbiased method such as exact diagonalization on the $`N`$-site lattice. Here we set $`N`$=8, as shown in Fig. 6(a). The correlation function of the multipole operators is given by
$$\chi _𝐪^{\mathrm{\Gamma }_\gamma }=(1/N)\underset{𝐫,𝐫^{}}{}e^{i𝐪(𝐫𝐫^{})}X_𝐫^{\mathrm{\Gamma }_\gamma }X_𝐫^{}^{\mathrm{\Gamma }_\gamma },$$
(29)
where $`\mathrm{}`$ denotes the expectation value using the ground-state wave-function.
In Fig. 6 (b), we show results for the correlation functions. The interaction between $`\mathrm{\Gamma }_{2u}`$ moments ($`b_{10}`$) is large, but the correlation function of the $`\mathrm{\Gamma }_{2u}`$ moment is not enhanced, indicating that the frustration effect is significant for an Ising-like moment such as $`\mathrm{\Gamma }_{2u}`$. We find large values of correlation functions for $`J_z^{4u2}`$, $`T_z^{5u}`$, and $`O_{xy}`$ moments at $`𝐪`$=$`(0,0,1)`$. However, there is no term in the effective model which stabilizes $`O_{xy}`$ quadrupole order at $`𝐪`$=$`(0,0,1)`$. We note that either of $`\mathrm{\Gamma }_{4u2}`$ and $`\mathrm{\Gamma }_{5u}`$ ordered states can accompany $`\mathrm{\Gamma }_{5g}`$ quadrupole moments. Thus, the enhancement of $`O_{xy}`$ correlation function indicates an induced quadrupole moment in $`\mathrm{\Gamma }_{4u2}`$ or $`\mathrm{\Gamma }_{5u}`$ moment ordered states. Namely, the relevant interactions are $`b_2^{(2)}`$ and $`b_8`$, which stabilize the $`J_z^{4u2}`$ and $`T_z^{5u}`$ order, respectively, at $`𝐪`$=$`(0,0,1)`$.
Next we study the ordered state by applying mean-field theory to the simplified model including only $`b_2^{(2)}`$ and $`b_8`$. Since the coupling constant $`b_8`$ is slightly larger than $`b_2^{(2)}`$, $`\mathrm{\Gamma }_{5u}`$ ordered state should has lower energy than $`\mathrm{\Gamma }_{4u2}`$ ordered state. The interaction $`b_8`$ stabilizes longitudinal ordering of the $`\mathrm{\Gamma }_{5u}`$ moments, i.e., $`𝐓_𝐫^{5u}𝐪`$.
However, we cannot conclude that the ground state is the single-$`𝐪`$ state $`(T_{x𝐫}^{5u},T_{y𝐫}^{5u},T_{z𝐫}^{5u})(0,0,\mathrm{exp}[i2\pi z/a])`$, since there is a possibility of multi-$`𝐪`$ structures. For isotropic moments, single-$`𝐪`$ and multi-$`𝐪`$ structures have the same energy, and thus, anisotropy in the moment is important to determine the stable structure. Indeed, the $`\mathrm{\Gamma }_{5u}`$ moment has an easy axis along in the $`\mathrm{\Gamma }_8`$ subspace. Kubo ; Kiss In this case, we find that a triple-$`𝐪`$ state is most stable among the single-$`𝐪`$ and multi-$`𝐪`$ states, since it gains interaction energy in all the directions.
In fact, the mean-field ground-state of the simplified model is the longitudinal triple-$`𝐪`$ $`\mathrm{\Gamma }_{5u}`$ octupole state with four sublattices, i.e.,
$`T_{x𝐫}^{5u}`$ $`\mathrm{exp}[i2\pi x/a],`$ (30a)
$`T_{y𝐫}^{5u}`$ $`\mathrm{exp}[i2\pi y/a],`$ (30b)
$`T_{z𝐫}^{5u}`$ $`\mathrm{exp}[i2\pi z/a].`$ (30c)
This state accompanies the triple-$`𝐪`$ quadrupole moment Paixao
$`O_{yz𝐫}`$ $`T_{x𝐫}^{5u},`$ (31a)
$`O_{zx𝐫}`$ $`T_{y𝐫}^{5u},`$ (31b)
$`O_{xy𝐫}`$ $`T_{z𝐫}^{5u}.`$ (31c)
In Fig. 7, we show symmetry of the charge distribution with spin density in the triple-$`𝐪`$ $`\mathrm{\Gamma }_{5u}`$ octupole state. Note that this triple-$`𝐪`$ structure does not have frustration even in the fcc lattice. The ground state energy is $`4b_8`$ per site, and the transition temperature is given by $`k_\text{B}T_{5u}`$=$`4b_8`$. We also note that this triple-$`𝐪`$ $`\mathrm{\Gamma }_{5u}`$ octupole state has been proposed for NpO<sub>2</sub> phenomenologically. Paixao
Let us now evaluate physical quantities in the mean-field theory. Figures 8(a) and (b) show an $`H`$-$`T`$ phase diagram and the magnetic field dependence of the magnetization at $`T`$=0, respectively. Note that the magnetization is isotropic as $`H`$$``$0 due to the cubic symmetry. The bend for $`𝐇[001]`$ and the dip for $`𝐇[110]`$ in magnetization indicate transitions to the two-sublattice structures. There is anomaly in magnetization also for $`𝐇[111]`$ at the transition to the different sublattice structure, but it is very weak.
Under a high magnetic field, sublattice structures change, as shown in Fig. 9: For $`𝐇[001]`$, we obtain a two-sublattice structure with
$`T_{x𝐫}^{5u}`$ $`=0,`$ (32a)
$`T_{y𝐫}^{5u}`$ $`=0,`$ (32b)
$`T_{z𝐫}^{5u}`$ $`\mathrm{exp}[i2\pi z/a].`$ (32c)
For $`𝐇[110]`$, there appears a two-sublattice structure with
$`T_{x𝐫}^{5u}`$ $`0,`$ (33a)
$`T_{y𝐫}^{5u}`$ $`0,`$ (33b)
$`T_{z𝐫}^{5u}`$ $`=0.`$ (33c)
Finally, for $`𝐇[111]`$, we observe
$`T_{x𝐫}^{5u}`$ $`\mathrm{sin}[2\pi (yz)/a],`$ (34a)
$`T_{y𝐫}^{5u}`$ $`\mathrm{sin}[2\pi (zx)/a],`$ (34b)
$`T_{z𝐫}^{5u}`$ $`\mathrm{sin}[2\pi (xy)/a].`$ (34c)
Note also that the triple-$`𝐪`$ state is fragile under $`𝐇[110]`$: $`T_{z𝐫}^{5u}`$=0 with a four-sublattice structure for $`g_J\mu _\text{B}H/(k_\text{B}T_{5u})0.11`$ at $`T`$=0 \[this phase boundary is not shown in Fig. 8(a)\].
Figures 10(a) and 10(b) show the temperature dependence of the specific heat and magnetic susceptibility, respectively. At $`T`$=$`T_{5u}`$, there appear the specific heat jump and a cusp in the magnetic susceptibility. In contrast to the sc and bcc lattices, there occurs single phase transition at zero magnetic field in the case of the fcc lattice. Note also that the cusp structure in the magnetic susceptibility is rather strong compared with experimental results.Kubo:NpO2 ; Ross Such a quantitative disagreement with experiments is considered to originate from the suppression of $`\mathrm{\Gamma }_7`$ orbital in our model. The analysis of the $`j`$=5/2 sextet model on the fcc lattice is one of future problems.
## V Discussion and summary
We have constructed $`\mathrm{\Gamma }_8`$ models with hopping integrals through $`(ff\sigma )`$ bonding based on the $`j`$-$`j`$ coupling scheme. In order to study multipole ordering, we have derived an effective model by using the second-order perturbation theory with respect to $`f`$-$`f`$ hopping. By applying mean-field theory, we find different multipole ordered states depending on the lattice structure. For the sc lattice, a $`\mathrm{\Gamma }_{3g}`$ antiferro-quadrupole transition occurs at a finite temperature. As lowering temperature further, we find a ferromagnetic transition. For the bcc lattice, a $`\mathrm{\Gamma }_{2u}`$ antiferro-octupole ordering occurs first, and a ferromagnetic transition follows it. Finally, for the fcc lattice, with careful analysis, we conclude the appearance of the single phase transition to the triple-$`𝐪`$ $`\mathrm{\Gamma }_{5u}`$ octupole ordering.
In this paper, we have not taken into account the effect of conduction electron. One may complain about this point, since it is believed that the hybridization of $`f`$ electrons with conduction electron band is important to understand the magnetism of $`f`$-electron systems. In fact, in the traditional prescription, first we derive the Coqblin-Schrieffer model from the periodic Anderson model by evaluating the $`c`$-$`f`$ exchange interaction $`J_{\mathrm{cf}}`$ within the second-order perturbation in terms of the hybridization between $`f`$\- and conduction electrons.Coqblin Then, we derive the RKKY interactions again using the second-order perturbation theory with respect to $`J_{\mathrm{cf}}`$.
In general, the RKKY interactions are orbital dependent and interpreted as multipole interactions. Such orbital dependence originates from that of the hybridization. Note that the hybridization should occur only between $`f`$\- and conduction band with the same symmetry. Here we emphasize that the symmetry of $`f`$-electron state is correctly included in our calculations. Thus, the structure in the multipole interactions will not be changed so much, even if we consider the effect of hybridization with conduction band, as long as we consider correctly the symmetry of $`f`$ electron states.
Let us show an example to support our belief. Concerning the octupole ordering in NpO<sub>2</sub>, we have extended the present theory by further including the effect of $`p`$ electrons of oxygen anions.Kubo:fp Namely, we have constructed the so-called $`f`$-$`p`$ model, given in the form of
$$=_\mathrm{f}+_\mathrm{p}+_{\mathrm{hyb}},$$
(35)
where $`_\mathrm{f}`$ and $`_\mathrm{p}`$ denote the local $`f`$\- and $`p`$-electron terms, respectively, and $`_{\mathrm{hyb}}`$ is the hybridization between $`p`$\- and $`f`$-electrons through $`(pf\sigma )`$ and $`(pf\pi )`$. Then, it has been found that the structure in the multipole interactions of the effective model derived from the $`f`$-$`p`$ model is qualitatively the same as those obtained in the $`\mathrm{\Gamma }_8`$ model on the fcc lattice. In fact, we have found a finite parameter region of $`\mathrm{\Gamma }_{5u}`$ antiferro-octupole phase. Namely, the $`f`$-$`p`$ model on the fcc lattice has a tendency toward $`\mathrm{\Gamma }_{5u}`$ antiferro-octupole ordering, which has been already captured in the simple $`(ff\sigma )`$ model. This result suggests that the structure in multipole interactions is determined mainly by the symmetry of $`f`$-electron state. Most of the effect of hybridization can be included by changing effectively $`(ff\sigma )`$ in the multipole interactions shown in the present paper.
However, if the itinerant nature of $`f`$ electrons is increased due to the large hybridization and metallicity of the ground state becomes significant, the present approximation inevitably loses the validity and the effect of the conduction band should be important. In such a case, it is necessary to develop a theory on the basis of the orbital-degenerate periodic Anderson model in order to include the multipole fluctuations. It is one of future tasks.
## Acknowledgments
We thank M. Yoshizawa and H. Fukazawa for sending us preprints on SmRu<sub>4</sub>P<sub>12</sub> prior to publication. We also thank H. Harima, S. Kambe, N. Metoki, Y. Tokunaga, K. Ueda, R. E. Walstedt, and H. Yasuoka for useful discussions. One of the authors (K. K.) is grateful to H. Onishi for useful comments on numerical diagonalization. K. K. is supported by the REIMEI Research Resources of Japan Atomic Energy Research Institute. Another author (T. H.) is supported by a Grants-in-Aid for Scientific Research in Priority Area “Skutterudites” under the contract No. 16037217 from the Ministry of Education, Culture, Sports, Science, and Technology of Japan. T. H. is also supported by a Grant-in-Aid for Scientific Research (C)(2) under the contract No. 50211496 from Japan Society for the Promotion of Science. |
warning/0506/hep-ph0506279.html | ar5iv | text | # From dense-dilute duality to self duality in high energy evolution
## Abstract
I describe recent work on inclusion of Pomeron loops in the high energy evolution. In particular I show that the complete eikonal high energy evolution kernel must be selfdual.
###### Keywords:
Pomeron loops, high energy evolution
Last year has seen renewed attempts to understand Pomeron loop contributions to the high energy evolution of hadronic cross sections in QCD. In recent years the study of the high energy scattering has centered around the so called JIMWLK evolution equation balitsky ; JIMWLK ; cgc . It describes the approach of the scattering amplitude to saturation due to multiple scattering corrections on dense hadronic targets, or in the diagrammatic language, the fan diagrams. The JIMWLK equation however only partially takes into account the processes whereby the gluons emitted in the projectile wave function at an early stage of the evolution, are ”bleached” by subsequently emitted gluons, or the so-called Pomeron loopsim1 ,ms ,ploops ,kl ,kl2 .
Recently we have calculated corrections to the JIMWLK equation, which take into account some finite density effects in the projectile wave function (or equivalently, resum certain corrections away form the dense limit of the the target)kl1 . We have also derived the evolution equation valid for dilute target, which is the opposite limit to that considered in JIMWLKkl . The most striking feature of the two results, is that they appear to be dual to each other. The improved JIMWLK equation is given bykl1
$`\chi ^{\mathrm{JIMWLK}+}`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle _z}\{b_i^a(z,0,[{\displaystyle \frac{\delta }{\delta \alpha }}])b_i^a(z,0,[{\displaystyle \frac{\delta }{\delta \alpha }}])+b_i^a(z,1,[{\displaystyle \frac{\delta }{\delta \alpha }}])b_i^a(z,1,[{\displaystyle \frac{\delta }{\delta \alpha }}])`$ (1)
$``$ $`2b_i^a(z,0,[{\displaystyle \frac{\delta }{\delta \alpha }}])\left[𝒫e^{i_0^1𝑑x^{}T^c\alpha ^c(x^{},z)}\right]^{ab}b_i^b(z,1,[{\displaystyle \frac{\delta }{\delta \alpha }}])\}`$
where $`𝒫`$ denotes path ordering with respect to $`x^{}`$ and the field $`b_i^a`$ satisfies the ”classical” equation of motionkl1 . The low density limit evolution kernel (KLWMIJ) including the same type of corrections but in the target wave function derived in kl is<sup>1</sup><sup>1</sup>1Following the work kl , the same expression has also been obtained in smith using the effective action techniques.:
$`\chi ^{\mathrm{KLWMIJ}+}=`$ $``$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle _z}\{b_i^a(z,0,[\rho ])b_i^a(z,0,[\rho ])+b^a(z,1,[\rho ])b^a(z,1,[\rho ])`$ (2)
$``$ $`2b_i^a(z,0,[\rho ])\left[𝒫e^{_0^1𝑑x^{}T^c\frac{\delta }{\delta \rho ^c(x^{},z)}}\right]^{ab}b_i^b(z,1,[\rho ])\}.`$
The two kernels are strikingly similar which suggests an intriguing duality between the high and the low density limits of the evolution kernel.
In this contribution I follow kl2 and show that indeed the full eikonal kernel for the high energy evolution must satisfy the property of self duality. The requirement that the evolution of the projectile and the target wave functions has the same functional form coupled with the requirement of Lorentz invariance of the scattering matrix, leads to the condition that the kernel of the evolution $`\chi [\rho ,\frac{\delta }{\delta \rho }]`$ must satisfy
$$\chi [\alpha ,\frac{\delta }{\delta \alpha }]=\chi [i\frac{\delta }{\delta \rho },i\rho ].$$
(3)
where $`\rho `$ is the charge density in the target wave function and $`\alpha `$ is defined byJIMWLK ; cgc
$`\alpha ^a(x,x^{})T^a={\displaystyle \frac{1}{^2}}(xy)\left\{U^{}(y,x^{})\rho ^a(y,x^{})T^aU(y,x^{})\right\},`$
$`U(x,x^{})=𝒫\mathrm{exp}\{i{\displaystyle _{\mathrm{}}^x^{}}𝑑y^{}T^a\alpha ^a(x,y^{})\}.`$ (4)
I note that from the functional integral point of view this duality has been discussed earlier in dual .
Consider the general expression for the $`S`$-matrix of a projectile with the wave function $`|P`$ scattering on a target with the wave function $`|T`$kl2 , where the total rapidity of the process is $`Y`$. The projectile is assumed to be moving to the left with rapidity $`YY_0`$ (and thus has sizeable color charge density $`\rho ^{}`$), while the target is moving to the right with rapidity $`Y_0`$ (and has large $`\rho ^+`$). We assume that the projectile and the target contain only partons with large $`k^{}`$ and $`k^+`$ momenta respectively: $`k^{}>\mathrm{\Lambda }^{}`$ and $`k^+>\mathrm{\Lambda }^+`$. The eikonal expression for the $`S`$-matrix reads
$$𝒮_Y=D\rho ^{+a}(x,x^{})W_{Y_0}^T[\rho ^+(x^{},x)]\mathrm{\Sigma }_{YY_0}^P[\alpha ],$$
(5)
where $`\mathrm{\Sigma }^P`$ is the $`S`$-matrix averaged over the projectile wave function
$$\mathrm{\Sigma }^P[\alpha ]=P|𝒫e^{i{\scriptscriptstyle 𝑑x^{}d^2x\widehat{\rho }^a(x)\alpha ^a(x,x^{})}}|P.$$
(6)
where $`W^T[\alpha ]`$ is the weight function representing the target, which is related to the target wave function in the following way: for an arbitrary operator $`\widehat{O}[\widehat{\rho }^+]`$
$$T|\widehat{O}[\widehat{\rho }^+(x)]|T=D\rho ^{+a}W^T[\rho ^+(x^{},x)]O[\rho ^+(x,x^{})].$$
(7)
The field $`\alpha (x)`$ is the $`A^+`$ component of the vector potential in the light cone gauge $`A^{}=0`$. This is the natural gauge from the point of view of partonic interpretation of the projectile wave function. In the formulae above we use hats to denote quantum operators. Note that the quantum operators $`\widehat{\rho }^a(x)`$ and $`\widehat{\rho }^{+a}(x)`$ do not depend on longitudinal coordinates, but only on transverse coordinates $`x`$. The ”classical” variables $`\alpha `$ and $`\rho ^+`$ on the other hand do depend on the longitudinal coordinate $`x^{}`$. This dependence, as discussed in detail in kl arises due to the need to take correctly into account the proper ordering of noncommuting quantum operators. Thus the ordering of the quantum operators $`\widehat{\rho }^+`$ in the expansion of $`\widehat{O}`$ in the lhs of eq.(7) translates into the same ordering with respect to the longitudinal coordinate $`x^{}`$ of $`\rho ^+(x^{})`$ in the expansion of $`O[\rho ^+(x^{})]`$ in the rhs of eq.(7).
As shown in kl the functional $`W^T[\alpha ]`$ cannot in general be interpreted as probability density, as it contains a complex factor. This factor - the Wess-Zumino term, ensures correct commutators between the quantum operators $`\widehat{\rho }^a`$. In the present derivation we do not require an explicit form of this term, but the following property which is implicit in eq. (7) is crucial to our discussion. The ”correlators” of the charge density $`\rho ^{a_1}(x_1,x_1^{})\mathrm{}\rho ^{a_n}(x_n,x_n^{})`$ do not depend on the values of the longitudinal coordinates $`x_i^{}`$, but only on their orderingkl .
Note that one can define an analog of $`W^T`$ for the wave function of the projectile via
$$P|\widehat{O}[\widehat{\rho }^{}(x)]|P=D\rho ^aW^P[\rho ^{}]O[\rho ^{}(x,x^{})].$$
(8)
With this definition it is straightforward to see that $`\mathrm{\Sigma }^P`$ and $`W^P`$ are related through a functional Fourier transform. To represent $`\mathrm{\Sigma }`$ as a functional integral with weight $`W^P`$ we have to order the factors of the charge density $`\widehat{\rho }^{}`$ in the expansion of eq.(6), and then endow the charge density $`\widehat{\rho }^{}(x)`$ with an additional coordinate $`t`$ to turn it into a classical variable. This task is made easy by the fact that the ordering of $`\widehat{\rho }`$ in eq.(6) follows automatically the ordering of the coordinate $`x^{}`$ in the path ordered exponential. Since the correlators of $`\rho (x,t_i)`$ with the weight $`W^P`$ depend only on the ordering of the coordinates $`t_i`$ and not their values, we can simply set $`t=x^{}`$. Once we have turned the quantum operators $`\widehat{\rho }`$ into the classical variables $`\rho (x^{})`$, the path ordering plays no role anymore, and we thus have
$$\mathrm{\Sigma }^P(\alpha )=D\rho ^aW^P[\rho ]e^{i{\scriptscriptstyle 𝑑x^{}d^2x\rho ^a(x,x^{})\alpha ^a(x,x^{})}}.$$
(9)
We now turn to the discussion of the evolution. The evolution to higher energy can be achieved by boosting either the projectile or the target. The resulting $`S`$-matrix should be the same. This is required by the Lorentz invariance of the $`S`$-matrix. Consider first boosting the projectile by a small rapidity $`\delta Y`$. This transformation leads to the change of the projectile $`S`$-matrix $`\mathrm{\Sigma }`$ of the form
$$\frac{}{Y}\mathrm{\Sigma }^P=\chi ^{}[\alpha ,\frac{\delta }{\delta \alpha }]\mathrm{\Sigma }^P[\alpha ]$$
(10)
Substituting eq.(10) into eq.(5) we have
$`{\displaystyle \frac{}{Y}}𝒮_Y`$ $`=`$ $`{\displaystyle D\rho ^{+a}(x,x^{})W_{Y_0}^T[\rho ^+(x^{},x)]\left\{\chi ^{}[\alpha ,\frac{\delta }{\delta \alpha }]\mathrm{\Sigma }_{YY_0}^P[\alpha ]\right\}}`$ (11)
$`=`$ $`{\displaystyle D\rho ^{+a}(x,x^{})\left\{\chi [\alpha ,\frac{\delta }{\delta \alpha }]W_{Y_0}^T[\rho ^+(x^{},x)]\right\}\mathrm{\Sigma }_{YY_0}^P[\alpha ]}.`$
Where the second equality follows by integration by parts. We now impose the requirement that the $`S`$-matrix does not depend on $`Y_0`$ LL . Since $`\mathrm{\Sigma }`$ in eq.(5) depends on the difference of rapidities, requiring that $`𝒮/Y_0=\mathrm{\hspace{0.17em}0}`$ we find that $`W`$ should satisfy
$$\frac{}{Y}W^T=\chi [\alpha ,\frac{\delta }{\delta \alpha }]W^T[\rho ^+]$$
(12)
Thus we have determined the evolution of the target eq.(12) by boosting the projectile and requiring Lorentz invariance of the $`S`$-matrix. On the other hand the extra energy due to boost can be deposited in the target rather than in the projectile. How does $`W^T`$ change under boost of the target wave function? To answer this question we consider the relation between $`\mathrm{\Sigma }`$ and $`W`$ together with the evolution of $`\mathrm{\Sigma }`$. Referring to eqs.(9) and (10) it is obvious that multiplication of $`\mathrm{\Sigma }^P`$ by $`\alpha `$ is equivalent to acting on $`W^P`$ by the operator $`i\delta /\delta \rho `$, and acting on $`\mathrm{\Sigma }^P`$ by $`\delta /\delta \alpha `$ is equivalent to multiplying $`W^P`$ by $`i\rho `$. Additionally, the action of $`i\rho `$ and $`i\delta /\delta \rho `$ on $`W^P`$ must be in the reverse order to the action of $`\delta /\delta \alpha `$ and $`\alpha `$ on $`\mathrm{\Sigma }^P`$. This means that the evolution of the functional $`W^P`$ is given by
$$\frac{}{Y}W^P=\chi [i\frac{\delta }{\delta \rho },i\rho ]W^P[\rho ].$$
(13)
Although eq.(13) refers to the weight functional representing the projectile wave function, clearly the functional form of the evolution must be the same for $`W^T`$. Comparing eq.(12) and eq.(13) we find that the high energy evolution kernel must, as advertised, satisfy the selfduality relation eq.(3). This is the main result.
The selfduality of the kernel is somewhat similar (although different in detail) to the duality symmetry of a harmonic oscillator Hamiltonian $`px,xp`$. One thus hopes that it may eventually be of help in solving the complete evolution equation, once it is derived. |
warning/0506/cond-mat0506471.html | ar5iv | text | # Silver transport in GexSe1-x:Ag materials: ab initio simulation of a solid electrolyte
## I Introduction
The chalcogenide glasses are preferred semiconducting materials for applications. They have well-defined niches in fiber opticsAggarwal , optical recordingMitkova , phase change memoryOvshinsky , and other technologies. Ge-Se glasses have been particularly studied because of their ready glass formation, easy synthesis requirements, and good chemical stability. The basic structural units are Se chains and Ge-Se tetrahedra which may combine in a variety of ways. Defects (including homopolar bonds) exist in these systems as seen in both experiment and theory.
The GeSe system was the first in which formation of an intermediate phase was demonstrated experimentally by Boolchand et alBoolchand and further developed theoretically by Thorpe, Jacobs, Chubnysky and PhillipsThorpe . In binary Ge<sub>x</sub>Se<sub>1-x</sub> glasses, the self-organized phase exists in the $`0.20<x<0.254`$ range, with glasses at $`x<0.20`$ regarded as floppy while those with $`x>0.26`$ stressed rigid. Dynamic calorimetry measurements on the intermediate phase have led to the conclusion that such materials do not ageBoolchand 2002 , a feature that may be of importance in application of these materials.
Silver added to chalcogenide glass hosts has attracted widespread interest in soft condensed matter scienceMitkova 1999 ; Wang . The interest emerges in part from the extensive bulk glass forming tendency in the Ge-Se-Ag ternary, the spectacular enhancement (eight orders of magnitude) in electrical conductivity of glasses with Ag relative to the glassy chalcogenide hosts, and from light-induced effects such as photo-doping, photo-diffusion and photo-deposition. Although the mobile ions in amorphous materials have been studiedAngell their detailed dynamics in amorphous hosts still constitutes one of the unsolved problems of solid state ionics. The structure of Ge-Se-Ag glass has been investigated using several experimental methods, including X-ray diffractionFischer ; piarristeguy neutron diffraction with isotopic substitutionLee , EXAFSOldale , differential anomalous X-ray scattering (DAS)Dejus ; Dejus 1992 ; Westwood and Modulated Differential Scanning Calorimetry (MDSC) and Raman spectroscopyMitkova 1999 ; Wang . Despite this impressive database, the structure of the ternary Ge-Se-Ag glasses has not yet been completely determined. There continues to be a debate on basic aspects of the glass structure (i.e. homogeneity and Ag coordination) especially for Se rich glasses with more than 67% Se.
Experimental evidence for macroscopic phase separation in these materials has come from MDSC results, which indicate bimodal glass transition temperaturesWang . In these experiments, one T<sub>g</sub> is independent of glass composition, and identified with a Ag<sub>2</sub>Se glass phasepbnature , while the second T<sub>g</sub> that varies with glass composition is related to the Ge-Se backbone. For the time scales (and possibly also lengths scales) of our simulations, these effects do not emerge, but are an interesting challenge for the future.
In a recent Lettertafdad05 , we briefly reported the motion of Ag ions in glassy chalcogenide hosts, and demonstrated the existence of ion trapping centers, which are important for relaxation processes in disordered systemsjcp . Here, we provide detailed information about structural and electronic properties, and also give new information about temperature dependence of the trapping, and also the geometry of the traps. In this paper, we have focused upon Ag-doped glasses containing Ge 25 at.% and Se 75 at.% (we later call this GeSe<sub>3</sub>). This composition is near the intermediate phase (slightly into the stressed-rigid phase). To our knowledge, this is the first ab initio simulation of these materials.
The rest of this paper is organized as follows. In Section II, we describe the simulation procedure to fabricate the atomistic models, discuss the ab initio total energy functional and force code used and other approximations. In Section III, we describe the structural properties using conventional measures such as static structure factors, and also apply a novel wavelet method to explore intermediate range order. Electronic properties are briefly discussed in Section IV, and Section V is concerned with the dynamics of the Ag<sup>+</sup> ions in the amorphous matrix.
## II Model Generation
### II.1 Energy Functional and Interatomic Forces
For the simulations reported in this paper, we use FIREBALL2000 developed by Lewis and coworkersLewis . This code is an approximate ab initio density functional approach to the electronic structure, total energies and forces based upon pseudopotentials and a real-space local basis of slighty-excited pseudoatomic orbitals to represent the Kohn-Sham functions. The method uses separable pseudopotentials, and allows the use of double-zeta numerical basis sets and polarization orbitals. The calculation is undertaken entirely in real space, which provides substantial computational efficiency. The exchange-correlation energy was treated within the LDA, using the results of Ceperley and Alderceperley , as interpolated by Perdew and Zungerperdew (more intricate gradient corrected functionals are available if needed). The pseudopotential and pseudoatomic wave functions were generated in the Troullier-Martins formtroullier employing the scheme of Fuchs and Schefferfuchs . Hamiltonian and overlap matrix elements are precalculated on a numerical grid and the specific values needed for a particular instantaneous conformation are extracted from the tabulated values via interpolation. Naturally, the integral tables need to be generated only once, for a given set of atomic species, rather than performing quadratures “on the fly” during a MD run.
### II.2 Model Formation
The models described here were generated using the melt quenching method. We began by randomly placing atoms in a cubic supercell according to the desired composition \[for (GeSe<sub>3</sub>)<sub>0.90</sub>Ag<sub>0.10</sub> 54 germanium atoms, 162 selenium atoms and 24 silver atoms; for (GeSe<sub>3</sub>)<sub>0.85</sub>Ag<sub>0.15</sub> 51 germanium atoms, 153 selenium atoms and 36 silver atoms\] with the minimum acceptable distance between atoms 2 Å. The size of the cubic cells was chosen to make the density of these glasses close to experimental data. The box size of the 240 atom supercell of (GeSe<sub>3</sub>)<sub>0.90</sub>Ag<sub>0.10</sub> and (GeSe<sub>3</sub>)<sub>0.85</sub>Ag<sub>0.15</sub> are respectively 18.601 Åand 18.656 Åwith corresponding densitypiarristeguy 4.98 g/cm<sup>3</sup> and 5.03 g/cm<sup>3</sup>. The structures were annealed and we obtained well thermalized melts at 4800K. We took three steps to cool the cells. First, the cells were equilibrated to 1100 K for 3 ps; then they were slowly cooled to 300 K over approximately 5 ps. The MD time step was 2.5fs. Simple velocity rescaling was used for the dissipative dynamics. In the final step, the cells were steepest descent quenched to 0K and maximum forces smaller in magnitude than 0.02 eV/Å. All calculations were performed at constant volume using the $`\mathrm{\Gamma }`$ point to sample the Brillouin zone to compute energies and forces.
The use of MD to quench a model liquid is natural and in principle completely general, since it superficially mimics the process of glass formation. It has been utilized for structural calculations for GeSe<sub>2</sub>Cappelletti , As<sub>2</sub>Se<sub>3</sub>Jun , Ge-Se-AgIyetomi and other systems (though there are sometimes hints that there is too much “liquid-like” character frozen into the resulting models). With accurate force calculations, simulated melt quenching has produced disappointing results for liquids such as GeSeRoon and ternary glasses such as As-Ge-Secomment . We believe that the success with (GeSe<sub>3</sub>)<sub>1-x</sub>Ag<sub>x</sub> ($`x`$=0.10, 0.15) glasses is connected to the fact we are in a weakly overconstrained glass forming part of the vibrational phase diagrammikejim . The “cook and quench” method is less effective in the highly overconstrained regime. It is likely that hybrid schemes mixing experimental information and ab intio simulation (such as the “ECMR” methodecmr ) would be most effective in this composition regime.
## III Structural Properties
### III.1 Short range order and defects
We define short range order (SRO) as a length scale not longer than the second nearest neighbor distance. Table 1 gives an overview of the SRO in g-(GeSe<sub>3</sub>)<sub>0.90</sub>Ag<sub>0.10</sub> and g-(GeSe<sub>3</sub>)<sub>0.85</sub>Ag<sub>0.15</sub>. The average bond length and average coordination number closely agree with the available datapiarristeguy . With increasing Ag concentration, the average coordination number increases to 2.8 in g-(GeSe<sub>3</sub>)<sub>0.85</sub>Ag<sub>0.15</sub> from 2.51 in g-GeSe<sub>3</sub>Petri . Also listed is the second nearest neighbor distance. Table 2 lists the average bonding distances of different possible bonds present in the model. Different values of the Ag-Ag bond distance have been proposed. Using the differential anomalous scattering (DAS), Westwood et alWestwood obtained a value of 3.35 $`\AA `$ for the Ag-Ag distance, a bit longer than our observation or other datapiarristeguy . Both models contain structural defects. In addition to the normally coordinated Ge<sub>4</sub> and Se<sub>2</sub>: Ge<sub>3</sub>, Se<sub>3</sub> and Se<sub>1</sub> are present in both models. Table 3 summarizes the statistical distribution of the main structural components. A look at the table shows that Ge-Se, Se-Se, Ag-Se and Ag-Ge correlations depend upon $`x`$ with additional Ag modifying the Ge-Se and Se-Se bonding. By integrating the partial pair correlation function g<sub>αβ</sub>(r) (Fig. 4), we estimate the average coordination number of Ag in the models. We found an average coordination of 2.0 and 2.9 in g-(GeSe<sub>3</sub>)<sub>0.90</sub>Ag<sub>0.10</sub> and g-(GeSe<sub>3</sub>)<sub>0.85</sub>Ag<sub>0.15</sub> respectively. We note that this is hard to estimate without ambiguity because of the lack of a well-defined deep minimum in $`g(r)`$ after the first peak especially for the $`x=0.15`$ model.
When discussing the coordination of Ag one can follow the arguments of KastnerKastner who was the first to describe the bonding of metal atoms in chalcogenide glasses. However, we have to consider the fact that in addition to the covalent bond that is expected to form with the chalcogens, Ag offers three empty $`sp`$ orbitals and is surrounded by the lone-pair electrons of its chalcogen neighbors. The latter offer the opportunity for the formation of up to three coordinate bonds. A coordinate bond is similar to a covalent bond and has similar strength but the bonding electrons are supplied by one bonding partner (the chalcogen atom)Cotton . As a result, the lowest energy-bonding configuration for Ag in chalcogenide glasses is an overall neutral complex with Ag positively charged with the negative charge located on neighboring chalcogen or chalcogensFritzsche . The coordination of Ag thereby can vary, but the expected value would be in average close to three - with one covalent bond and up to 2 coordinated bonds. The opportunity for four-fold coordination also exists as there is one more free $`sp`$ orbital at the Ag atom but evidently this does not satisfy the requirements of the lowest energy configuration and the electronegativity of the entire complex in the presence of other cations in the system so that the probability for this type of bonding is lower. Our simulations appear to be consistent with these chemical considerations, as all of the Ag is two-fold in the 10% model, with the addition of three-fold Ag in the 15% Ag model.
### III.2 Intermediate range order
Fig. 1 shows the calculated static structure factors for (GeSe<sub>3</sub>)<sub>0.90</sub>Ag<sub>0.10</sub> and (GeSe<sub>3</sub>)<sub>0.85</sub>Ag<sub>0.15</sub> and the comparison with the experimental data from ref. piarristeguy . As we did not include a priori information in the model formation process, the fact that the peak positions and spectral weight of S(Q) agree well with experimental data is encouraging. The third and fourth peaks are a result of the short range order in the models. The position of the second peak does not depend strongly on the Ag concentration (though its width does). This peak is located at Q$``$2.09 Å<sup>-1</sup> in both models. By contrast the third peak intensity decreases with increasing Ag. This peak is located at Q$``$3.49 Å<sup>-1</sup> in (GeSe<sub>3</sub>)<sub>0.90</sub>Ag<sub>0.10</sub> and at Q$``$3.41 Å<sup>-1</sup> in (GeSe<sub>3</sub>)<sub>0.85</sub>Ag<sub>0.15</sub>; whereas in GeSe<sub>3</sub> (x=0, concentration of Ag in (GeSe<sub>3</sub>)<sub>1-x</sub>Ag<sub>x</sub>) this peak appears nearpiarristeguy ; Petri 3.53 Å<sup>-1</sup>. In Fig. 2 the partial structure factors show that S<sub>GeGe</sub>(Q) and S<sub>GeSe</sub>(Q) tend to cancel each other.
Both models reveal a feature in S(Q) near 1.07 Å<sup>-1</sup>. This is a harbinger of a First Sharp Diffraction Peak (FSDP), that becomes explicit in more Ge-rich materials. Piarristeguy et alpiarristeguy show that this peak varies as a function of Ag content. As Ag concentration increases, the FSDP decreases due to a change of the IRO. Moreover, Ag disturbs the GeSe<sub>4/2</sub> network, and leads to the fragmentation of GeSe<sub>4/2</sub> tetrahedrons. From the partial structure factors it is apparent that the “proto-FSDP” has contributions from all of the partials. To elucidate the intermediate range order more quantitatively, we have used the wavelet-based methods of Harrop and coworkersHarrop . This is a promising scheme for interpreting structure factor data based upon continuous wavelet transformswaveletnote (CWT). In Fig. 3, we illustrate the results of the wavelet analysis extracted from the experimental datapiarristeguy . The most obvious feature is the extended range real-space correlations associated with the diffraction peak near 3.5Å<sup>-1</sup> for the 10% Ag model. This is connected to the narrowness of the peak for the 10% glass relative to the 15% material: in the latter case the correlations disappear by about 15Å<sup>-1</sup>, whereas the correlations extend to at least 25Åfor the 10% glass. A similar state of affairs also accrues for the peak near 5.5Å<sup>-1</sup> and for a similar reason. This work emphasizes that simple associations of the reciprocal of the peak position of the FSDP to real-space length scales is misleading, particularly if the peak is narrow as discussed by Uchino and coworkers in studies of silica glassHarrop .
In Fig. 4 we plot the partial pair correlation function of both models. The Ge-Se and Se-Se pairs provide the dominant contribution to the first shell of the pair correlation function g(r) whereas Ag-Se contribute to the second peak, Ag-Ag to the third and Se-Se (second nearest neighbors) to the fourth peak.
## IV Electronic Properties
Having studied structural properties, we now briefly analyze the electronic properties of our models. The electronic density of states (EDOS) of both models are calculated and analyzed by the inverse participation ratio (IPR), which we denote by $``$. The EDOS are obtained by summing suitably broadened Gaussians centered at each eigenvalue. The IPR
$$(E)=N\underset{n=1}{\overset{N}{}}q(n,E)^2$$
determines the spatial localization of electronic eigenvalues. Here $`N`$ is the number of atoms in the model and $`q(n,E)`$ is the Mulliken charge localized on atomic site $`n`$ in a certain eigenstate $`E`$. Hence, $``$ is a measure of the inverse number of sites involved in the state with energy E. For a uniformly extended state, the Mulliken charge contribution per site is uniform and $``$(E)=1/N. For an ideally localized state, only one atomic site contributes all the charge and $``$(E)=1. Therefore a larger value of $``$ means that the eigenstate is more localized in real space.
In Fig. 5 we report the EDOS and the species-projected density of states of (GeSe<sub>3</sub>)<sub>0.90</sub>Ag<sub>0.10</sub> and (GeSe<sub>3</sub>)<sub>0.85</sub>Ag<sub>0.15</sub> glasses (our electronic eigenvalues have been shifted in order to place the valence band edge eigenvalue at zero). It should be noted that the spectra of both models are similar and closely related to the EDOS of Ge<sub>x</sub>Se<sub>1-x</sub> ($`x>0.15`$). With the addition of Ag into g-GeSe<sub>3</sub>, an intense peak, due to the Ag 4$`d`$ electrons appears at about -3.47 eV as shown in Fig. 5. The valence band exhibits three features. The two lowest bands between -14.8 eV and -7.0 eV originate from the atomic 4$`s`$-like states of Ge and Se partially hybridized to form bonding states to Ag atoms. The next band lying between -7.0 and 0.0 eV contains $`p`$ like bonding states of Ge and Se and $`d`$ like bonding states of Ag. The peak in the topmost valence region is due to the lone-pair 4$`p`$ electrons of Se atoms. The $`\mathrm{\Gamma }`$ point optical gaps of g-(GeSe<sub>3</sub>)<sub>0.90</sub>Ag<sub>0.10</sub> and g-(GeSe<sub>3</sub>)<sub>0.85</sub>Ag<sub>0.15</sub> are respectively of the order of 1.20 and 1.26 eVgap . As the Ag content increases, the optical band gap increases. To our knowledge, experimental information about the EDOS of both systems is unavailable, so the curve in Fig. 5 is actually a prediction. This is an interesting contrast to the work of Simdyankin and coworkerssim , who show that for low concentrations of Cu, the gap decreases with addition of Cu in AsS and AsSe glasses. Care is needed in comparing these results since the hosts and transition metals are different, and our models have far higher metal content.
In order to connect localized eigenstates to particular topological/chemical regularities we plot in Fig. 6 the IPR in the band gap region. We found that the localization in the valence band is mainly due to Se atoms and in the conduction band to Ge atoms in both models. A close look at the localized states at the band edges shows that the localized states at the top of the valence band of g-(GeSe<sub>3</sub>)<sub>0.90</sub>Ag<sub>0.10</sub> are mostly associated with two- and three-fold coordinated Se atoms with homopolar bonds, whereas the localization at the conduction band edge arises from overcoordinated Se associated with homopolar bonds and four-fold coordinated Ge connected to Se atoms involved with Se-Se homopolar bonds. By contrast, the top of the valence band of g-(GeSe<sub>3</sub>)<sub>0.85</sub>Ag<sub>0.15</sub> is quite extended; the conduction band edge shows very few localized states due to the overcoordinated Ge and Se. This explains the results of KawasakiKawasaki showing dominance of the ionic conductivity related to Ag<sup>+</sup> ions at these particular compositions.
## V The Dynamics of Silver Ions
An outstanding feature of the materials is the high mobility of the silver in the complex host. The chemical explanation for the high diffusivity of Ag is its high quadrupolar deformabilityFritzsche . To “mine” information from our MD simulations, we begin by computing the mean-square displacement (MSD) functions for all of the atomic constituents:
$$r^2(t)_\alpha =\frac{1}{N_\alpha }\underset{i=1}{\overset{N_\alpha }{}}|𝐫_𝐢(𝐭)𝐫_𝐢(0)|^2,$$
(1)
where $``$ is an average over MD simulation time and the sums are over particular atomic species, $`\alpha `$. The MSD were calculated for Ge, Se and Ag ions for both models and are shown in Fig. 7. The MSD of Ag ions increases rapidly with time, whereas that of the Ge and Se show a very slight slope. We chose a temperature of 1000 K to illustrate the diffusion, and later discuss behavior at selected lower temperatures.
To explain the mechanism of diffusion of silver, we examine the trajectories of these particles. We obtained 2.5$`\times `$10<sup>4</sup> steps of time development, for a total time of 62.5 ps and a fixed temperature of 1000 K. Fig. 8 illustrates 2D projections of trajectories of the most and least mobile Ag atoms in the $`x=0.15`$ model. We notice that for short times, the MSD of the most mobile atoms increases due to the diffusive motion of Ag. At intermediate times, the atoms may be trapped in a cage formed by their neighbors, and at the longest times we can explore, they can escape such traps and diffuse again. Thus, our trajectories can largely be separated into vibration around stable trapping sites and hops between such sites. In both glasses, a fraction of Ag atoms move large distances (see Fig. 9). In (GeSe<sub>3</sub>)<sub>0.90</sub>Ag<sub>0.10</sub>, about 91.67 % of silver atoms move an average distance greater than 2.5 Åfor a time scale of 39 ps. Among them 25 % have an average displacement greater than 5 Å. By contrast only 4.2 % of Ag atoms move less than 2 Å. On the other hand, about 89 % of Ag atoms move on an average distance greater than 2.5 Åin (GeSe<sub>3</sub>)<sub>0.85</sub>Ag<sub>0.15</sub> for the same period of time. 36.1 % of those atoms have an average displacement greater than 5 Å. The most mobile Ag atoms move on an average distance of 8.2 Å. These numbers illustrate the high ionic mobility of Ag ions in these complex glasses and are suggestive for a significant contribution of correlated hops of Ag<sup>+</sup> in the diffusion process.
Based on these trajectories, thermal transport coefficients such as diffusion coefficients can be evaluated. Typically one uses either the Green-Kubo formula Chandler where the time-dependent velocity autocorrelation (VAC) is integrated or the Einstein relation Chandler is employed, and the MSD is differentiated with respect to time. Since transport coefficients are equilibrium properties, the system must be properly thermalized before the transport properties can be estimated. Our simulations show that at time t$`>`$4ps the systems are well equilibrated . The Einstein relation for self-diffusion reads:
$$|𝐫(𝐭)𝐫(\mathrm{𝟎})|^\mathrm{𝟐}=\mathrm{𝟔}𝐃𝐭+𝐂$$
(2)
where C is a constant, $`D`$ is the self-diffusion coefficient and $`|𝐫(𝐭)𝐫(\mathrm{𝟎})|^\mathrm{𝟐}`$ is the mean-square distance from initial position at time $`t`$, averaged over atoms of a given species. Direct simulation of the atomic trajectory and simple fitting yields C and D, and, in particular, estimates for the self-diffusion coefficient of Ag, $`D_{Ag}`$. The estimated values of $`D_{Ag}`$ as a function of temperature are listed in Table 4. These results are qualitatively reasonable when compared to the recent (room temperature) experiments of Ure$`\stackrel{~}{n}`$a et alUrena with the appropriate exponential activation factor included. The probability of correlated motion of Ag<sup>+</sup> will of course increase with increasing Ag concentration in agreement with decrease of the activation energy for conductivityGutenev .
The dynamics is sensitive to the temperature. We performed additional simulations at temperatures ranging from 640K to 1000K. In Fig. 10 we illustrate the hopping, which is qualitatively like the high temperature hopping, and the traps are very well defined. We note that even the most mobile Ag ions spend a substantial time in the traps, and appear to hop very efficiently (quickly) between traps.
To further study the trapping centersPhillips we also obtained estimates of trap sizes and trap lifetimes as a function of temperature. In our calculations, we consider only particles that experience more than one traps. To calculate the size of the traps we enclose the particle trajectories of each trap in a sphere of radius $`r_{tr}`$ centered on the average position of the particle in the trap. Then, we determine the displacement of the particle trajectories in the trap with respect to the average trajectory, and average over all the displacements. The trap size $`r_{tr}`$ is then obtained by averaging over different traps. Knowing the trap sizes, their lifetimes can be easily determined. In Table 4, we give an estimate of the trap sizes and trap lifetimes as a function of temperature. As T increases, the trap sizes increase and their lifetimes decrease. The averaged trap radii range from about 0.8 to more than 1.8 Å, and the averaged lifetimes extend from about 2.5 to more than 5 ps. The individual trap radii can be as large as 2.4 Åand the lifetimes as long as 7 ps. We have also used the MSD methodweitz to estimate trap sizes, and obtain results within a factor of $`2`$ from our simple geometrical approach.
Further insight into the mechanism of diffusion in (Ge<sub>x</sub>Se<sub>1-x</sub>)<sub>1-y</sub>Ag<sub>y</sub> can be found by studying the behavior of the molar volume of particular regions containing silver atoms. Hence we calculate the local density of the most and least mobile silver atoms as a function of time, then compare them to the density of the glass. To do so, we draw a sphere of radius R=4 Å. The center of the sphere is the position of the Ag atoms we are tracking at a time t (the center of the sphere varies as a function of time). Then we calculate the mean density of atoms inside the sphere. Fig. 11 illustrates the local density of a few Ag atoms as a function of time. As seen on the figure, the most mobile Ag atoms are consistently located in regions with a lower local density (lower local volume fraction) and higher disorder. On the other hand, as we showed by direct calculation, there is little if any correlation between the trajectory-averaged mean density and the tendency to diffuse for the Ag atoms. On the other hand, perhaps unsurprisingly, we found also explored correlations between the average displacement of mobile Ag ions and the standard deviation of their local density ($`\sigma _i`$=$`\sqrt{\rho _i^2\rho _i^2}`$, where $`\rho _i`$ is the local density of the Ag ion $`i`$, and $``$ means trajectory average.). We found that as the average displacement increases, the standard deviation becomes larger (see Fig. 12). The correlations are perhaps linear with much noise. This implies that the diffusive Ag ions are exploring a wide variety of densities and the weakly-diffusing Ag sample a restricted density range.
## VI Conclusion
We have presented ab initio models of GeSe glasses heavily doped with Ag and studied the dynamics of the network with an emphasis on the motion of Ag ions. The models reproduce structural data, including reasonably subtle features in the diffraction data including the first peak (or shoulder) in S(Q). Wavelet methods help significantly in revealing intermediate range real-space correlations in the glasses. The atomistic motion of Ag<sup>+</sup> ions is detailed for short times with a reliable first principles interaction. We have shown by direct calculation that trapping centers exist, and have shown that local basis ab initio MD can provide direct insight into the processes of transition metal dynamics in amorphous chalcogenide materials.
## VII Acknowledgments
We thank the US National Science Foundation for support under grants DMR-0074624, DMR-0205858 and DMR-0310933. We also gratefully acknowledge the support of Axon Technologies, Inc. We thank Dr. J. C. Phillips for many helpful comments and insights, and Drs. Jon Harrop, S. N. Taraskin and Professor S. R. Elliott (University of Cambridge) for sharing their wavelet techniques for the analysis of intermediate range order. |
warning/0506/cond-mat0506710.html | ar5iv | text | # No sliding in time
## 1 Introduction
When discussing the properties of a quantum many-body Hamiltonian, one commonly intermixes the notions of *gaplessness* and *criticality*. These two properties, however, may or may not automatically imply each other. The former refers to the absence of an energy gap separating the GS(s) from the excited states and implies slower than exponential, typically power-law decay of correlations in imaginary time. The latter refers to the equal-time correlations of operators separated in space. Naturally, if the system is Lorentz-invariant, the two properties are identical. Even in the absence of Lorentz invariance, these two properties typically follow from each other. *E.g.*, at a quantum critical point, there is a dynamical critical exponent $`z`$ relating the scaling of correlations in space and in time: $`tx^z`$. Since all standard examples of quantum critical points are characterised by finite non-zero values of $`z`$, gaplessness once again implies criticality and vise versa. However, a spin model introduced recently in the context of systems with topological order , was shown to be gapless while possessing only short-ranged equal time correlations between *local* operators. Moreover, in that particular model one could add another term to the Hamiltonian which would open a gap without affecting equal-time correlations. This last feature, of course, should not come as a surprise since for any model with short-ranged ground-state correlations, one could just add to the Hamiltonian an operator $`(𝕀𝒫_0)`$ where $`𝒫_0`$ is a projector onto this ground state. Such a projector will necessarily open a spectral gap while the quasi-local nature of ground state correlations will almost certainly lead to the quasi-local nature of the contributing terms in the Hamiltonian. It is the former feature, namely the spectral gaplessness of a system with short-range correlations that comes as a surprise <sup>1</sup><sup>1</sup>1Notice that such a possibility is remarkable from the point of view of understanding *glasses*: one of the puzzles about a generic glassy behaviour is their slow dynamics with no long spatial correlations.. As discussed in , this is due to a “bottle-neck” quantum dynamics that does not allow for the efficient mixing of different quantum states contributing to the ground state and thus allows one to construct a “twisted” excited state (in the spirit of Lieb-Schultz-Mattis theorem ) whose energy is vanishing in the thermodynamic limit. A similar situation was encountered in models where the “bottle-neck” quantum dynamics, in addition to leading to energy gaps that are exponentially small in domain sizes, prevents the systems from reaching the true ground state when coupled to a quantum bath at zero temperature, remaining in a state of “quantum glassiness” .
A natural question then arises: could we also find a system described by a *local* Hamiltonian with a spectral gap which nevertheless displays power-law decay of equal-time correlations of *local* operators. The condition of locality is important here: relaxing the requirement of locality of a Hamiltonian leads to a trivial “yes” by the means of adding a projector onto the critical ground state as described in the previous paragraph. We also know examples of gapped systems with quasi-long ranged correlations of non-local operators such as the non-local order parameters in the quantum Hall effect and other topological phases . In short, the answer to the above question is “No” due to the theorem proved by Hastings . This, however, contradicts our intuition drawn from Statistical Mechanics where *sliding* phases are known to exist, as has been demonstrated in the context of stacks of 2D layers of $`XY`$ spins with gradient coupling between the layers . The resulting *sliding* phase is characterised by an algebraic decay of in-plane correlations within the layer and an exponential decay in the perpendicular direction. Later, this construction has been extended to the quantum systems of stripes , with the layers now representing 1D stripes in space-time. Such phases were also found in dissipative Josephson junction arrays .
In general, it is common in the study of quantum critical phenomena of $`d`$-dimensional systems to relate the problem to $`(d+1)`$-dimensional classical systems. This can be done generically, and perhaps one place where a true difference between quantum systems in Euclidean time and classical systems may arise is when a topological term exists and the resulting classical action cannot be made real. Typically, in going from a $`d`$-dimensional quantum system with local interactions and dynamics to a $`(d+1)`$-dimensional classical system, one maps the local quantum Hamiltonian density into a local classical Lagrangian density. Hence, if the critical properties of this local classical system can be understood, so would those of the quantum model. Now, one could turn the question around, and ask instead about a behaviour of a model derived from a local classical Lagrangian. Here we would like to address this question in a particularly interesting case - that of a classical sliding phase. A quantum version of this model would imply critical correlations in space but not in time, if the direction perpendicular to the layers is taken to be Euclidean time. Thus the system is expected to be gapped and yet have critical correlations in space, seemingly a violation of Hastings’ theorem. The purpose of this Letter is to examine this apparent contradiction.
## 2 Sliding Phase
We now quickly review the sliding phase as described in . The *classical* Hamiltonian consists of three parts:
$$H=H_0+H_\mathrm{g}+H_\mathrm{J}.$$
(1)
Here $`H_0`$ is just a sum over all layers of independent $`XY`$-Hamiltonians:
$$H_0=\frac{K}{2}\underset{n}{}d^2r\left[\mathbf{}_{}\theta _n(𝐫)\right]^2,$$
(2)
where $`𝐫=(x,y,0)`$ is a point in the $`x`$-$`y`$ plane and $`\mathbf{}_{}`$ is the gradient operator acting on these two coordinates.
The second term in Eq. (1) couples gradients of $`\theta _n`$ in different layers:
$$H_\mathrm{g}=\frac{1}{2}\underset{n,m}{}d^2r\frac{U_m}{2}\left\{\mathbf{}_{}\left[\theta _{n+m}(𝐫)\theta _n(𝐫)\right]\right\}^2.$$
(3)
Finally, Josephson couplings between layers are added:
$$H_\mathrm{J}=V_\mathrm{J}d^2r\mathrm{cos}\left[\theta _{n+p}(𝐫)\theta _n(𝐫)\right].$$
(4)
Following , we choose to consider only two-layer couplings with $`p`$ being the distance between the coupled layers. For the reasons that will become clear shortly, we will concentrate on the next-nearest layer coupling, $`p=2`$. As follows from , the physics of the sliding phase should not depend on this choice.
In the absence of the gradient couplings (second term in Eq. (1)), this Hamiltonian is just that of a 3D $`XY`$-model (strictly speaking, for $`p=2`$, it is two independent interlaced $`XY`$-models). It might appear unorthodox that we keep the cosine coupling between the layers while using its expanded form for the in-plane couplings. The continuous in-plane limit is not an issue here; we simply adopt it from . This choice of couplings is permitted keeping in mind the ultimate goal of constructing the phase with critical intra-layer and exponential inter-layer correlations. In such a phase, one can ignore in-plane vortices but must allow for between-the-planes ones. However, without the gradient couplings, there is no such phase as the temperature at which layers decouple turns out to be higher than the Kosterlitz-Thouless (KT) temperature in a single layer.
On the other hand, in the absence of Josephson couplings ($`V_J=0`$), one obtains the *ideal* sliding Hamiltonian $`H_\mathrm{S}=H_0+H_\mathrm{g}`$, which is invariant with respect to $`\theta _n(𝐫)\theta _n(𝐫)+\psi _n`$ for any constant $`\psi _n`$. I.e., the energy is unchanged when angles in different layers slide relative to one another by arbitrary amounts – one can think of it as a “reduced” gauge symmetry with one gauge choice per layer. As a result, the angles in different layers are uncorrelated. The low-temperature phase is the *ideal* sliding phase characterised by $`\mathrm{cos}[\theta _m(𝐫)\theta _n(0)]\delta _{m,n}r^\eta `$ with $`\eta =T/(2\pi \stackrel{~}{K})`$, where $`\stackrel{~}{K}`$ is the renormalised in-plane coupling. For the simplest case of only nearest layer coupling $`U_m=U\delta _{m,\pm 1}`$, the effective in-plane coupling is $`\stackrel{~}{K}=K\sqrt{1+4U/K}`$.
What constitutes the main result of , is the fact that in the presence of *both* the gradient and Josephson couplings between the layers, a careful choice of coupling constants may open a “window of opportunity” in temperature within which the layers decouple before the spins in each layer completely disorder via a KT transition. In other words, the inter-layer vortices proliferate before the in-plane vortices unbind. This is the sliding phase which we shall now discuss in the quantum context.
## 3 Quantum Hamiltonian
Our discussion of constructing the corresponding quantum Hamiltonian will follow Kogut’s review . If the classical Hamiltonian $`H`$ consists of only intra-layer terms, $`H^{}[\theta _n]`$ and terms coupling neighbouring layers, $`H^{\prime \prime }[\theta _n,\theta _{n+1}]`$, the corresponding partition function can be written as a $`Z=\mathrm{tr}\left(\widehat{T}^N\right)`$ where the interlayer transfer matrix is defined as
$$\theta _{n+1}\left|\widehat{T}\right|\theta _n=\mathrm{exp}\left\{\beta \left(\frac{1}{2}H^{}[\theta _n]+\frac{1}{2}H^{}[\theta _{n+1}]+H^{\prime \prime }[\theta _n,\theta _{n+1}]\right)\right\}$$
(5)
and $`N`$ is the total number of layers. We now identify the $`n`$-th layer with the imaginary time slice $`\tau _n`$ while $`\tau _{n+1}=\tau _n+ϵ`$. The quantum Hamiltonian $`\widehat{}`$ is then defined by $`\theta (\tau _n+ϵ)\left|\mathrm{exp}\left\{ϵ\widehat{}\right\}\right|\theta (\tau _n)=\theta (\tau _n+ϵ)\left|\widehat{T}\right|\theta (\tau _n)`$ in the limit of $`ϵ0`$. Notice that the explicit form of a quantum Hamiltonian may be difficult to obtain: the transfer matrix needs not be in the form of an exponential of a simple operator expression. It is, however, often easy to represent it as a product of such exponentials, each corresponding to a term in the classical Hamiltonian.
### 3.1 Ideal sliding
Let us first look at the case of an ideal sliding Hamiltonian ($`V_J=0`$). For simplicity, for now we also restrict ourselves to the nearest-neighbour coupling: $`U_m=U\delta _{m,\pm 1}`$. Denoting $`\theta =\theta _n=\theta ^{(i)}(\tau )`$ and $`\theta ^{}=\theta _{n+1}=\theta ^{(i)}(\tau +ϵ)`$ we have:
$$\theta ^{}|\widehat{T}_\mathrm{S}|\theta =\mathrm{exp}\left[\frac{\beta }{4}d^2r\left(K\left[\mathbf{}\theta (𝐫)\right]^2+K\left[\mathbf{}\theta ^{}(𝐫)\right]^2+2U\left\{\mathbf{}\left[\theta ^{}(𝐫)\theta (𝐫)\right]\right\}^2\right)\right].$$
(6)
Introducing the canonically conjugate quantised fields $`\widehat{\theta }(𝐫)`$ and $`\widehat{\pi }(𝐫)`$ such that $`[\widehat{\pi }(𝐫),\widehat{\theta }(𝐫^{})]=i\delta (𝐫^{}𝐫)`$, we obtain:
$`\widehat{T}_\mathrm{S}`$ $``$ $`\mathrm{exp}\left\{{\displaystyle \frac{\beta K}{4}}{\displaystyle d^2r\left[\mathbf{}\widehat{\theta }(𝐫)\right]^2}\right\}`$ (7)
$`\times `$ $`\mathrm{exp}\left\{{\displaystyle \frac{1}{4\pi \beta U}}{\displaystyle d^2rd^2r^{}\widehat{\pi }(𝐫)\widehat{\pi }(𝐫^{})\mathrm{ln}\left|𝐫𝐫^{}\right|}\right\}`$
$`\times `$ $`\mathrm{exp}\left\{{\displaystyle \frac{\beta K}{4}}{\displaystyle d^2r\left[\mathbf{}\widehat{\theta }(𝐫)\right]^2}\right\}.`$
This can be explicitly verified by substituting the above expression for $`\widehat{T}`$ into Eq. (6). Notice that the operators in the first and the third exponent of Eq. (7) are diagonal in the $`\theta `$ representation while for the second exponential we have:
$`\theta ^{}|\mathrm{exp}\left\{g{\displaystyle d^2rd^2r^{}\widehat{\pi }(𝐫)\widehat{\pi }(𝐫^{})\mathrm{ln}\left|𝐫𝐫^{}\right|}\right\}|\theta `$
$`={\displaystyle 𝒟p𝒟p^{}\theta ^{}|p^{}p^{}|\mathrm{exp}\left\{gd^2rd^2r^{}\widehat{\pi }(𝐫)\widehat{\pi }(𝐫^{})\mathrm{ln}\left|𝐫𝐫^{}\right|\right\}|pp|\theta }`$
$`={\displaystyle 𝒟p\mathrm{exp}\left\{d^2r\left[\mathrm{i}p(𝐫)[\theta ^{}(𝐫)\theta (𝐫)]+gd^2r^{}p(𝐫)p(𝐫^{})\mathrm{ln}\left|𝐫𝐫^{}\right|\right]\right\}}`$
$`{\displaystyle 𝒟p_k\mathrm{exp}\left\{\frac{d^2k}{(2\pi )^2}\left[\mathrm{i}p_𝐤\left(\theta _𝐤^{}\theta _𝐤\right)\frac{2\pi g}{k^2}p_𝐤p_𝐤\right]\right\}}`$
$`\mathrm{exp}\left\{{\displaystyle \frac{1}{8\pi g}}{\displaystyle \frac{d^2k}{(2\pi )^2}k^2\left(\theta _𝐤^{}\theta _𝐤\right)\left(\theta _𝐤^{}\theta _𝐤\right)}\right\}`$
$`=\mathrm{exp}\left({\displaystyle \frac{1}{8\pi g}}{\displaystyle d^2r\left\{\mathbf{}\left[\theta ^{}(𝐫)\theta (𝐫)\right]\right\}^2}\right).`$ (8)
Instead of writing the corresponding quantum Hamiltonian formally as $`\widehat{}_\mathrm{S}=(1/ϵ)\mathrm{ln}\widehat{T}_\mathrm{S}`$, we can simplify it by going to the continuous time limit $`ϵ0`$. We introduce the new coupling constants, $`g_\theta \frac{\beta K}{2ϵ}`$ and $`g_\pi \frac{1}{4\pi \beta Uϵ}`$ and require that they do not scale with $`ϵ`$, which implies the following scaling for the original couplings: $`Kϵ`$, $`Uϵ^1`$. Then, up to the corrections of order $`ϵ`$,
$$\widehat{}_\mathrm{S}=g_\theta d^2r\left[\mathbf{}\widehat{\theta }(𝐫)\right]^2g_\pi d^2rd^2r^{}\widehat{\pi }(𝐫)\widehat{\pi }(𝐫^{})\mathrm{ln}\left|𝐫𝐫^{}\right|.$$
(9)
Notice that the implied scaling of $`K`$ and $`U`$ does not lead to any problems with tuning the parameters to the values needed to reach the desired sliding phase. This is because the *effective* coupling governing the behaviour of the classical statistical mechanical model is given by
$$\beta \stackrel{~}{K}=\beta K\sqrt{1+\frac{4U}{K}}=2ϵg_\theta \sqrt{1+\frac{1}{2\pi g_\theta g_\pi ϵ^2}}\sqrt{\frac{2g_\theta }{\pi g_\pi }}\mathrm{as}ϵ0.$$
(10)
Thus, this model has a well-defined continuous time limit.
However, the quantum version of the ideal sliding Hamiltonian contains logarithmically long ranged interactions between momenta at different points. As will be discussed below, Hastings’ theorem does not apply to such a Hamiltonian, hence no contradiction appears here. We remark on an interesting physical picture arising from the quantum Hamiltonian (9) if we think of the eigenvalues of $`\widehat{\pi }`$ as “charge” or “vorticity”. Notice that due to the compactness of the conjugate variable $`\theta `$, the eigenvalues of $`\widehat{\pi }`$ are quantised in integer units (we used $`\mathrm{}=1`$), this is completely analogous to the quantisation of $`L_z`$ in quantum mechanics. Therefore the second term in the quantum Hamiltonian (9) describes a classical 2D Coulomb gas (or a gas of vortices) with the usual logarithmic interaction. We know that the collective mode (plasmon) in such gas is gapped (unlike in the case of $`1/r`$ interactions). The first term has no simple classical meaning in this language; it is responsible for making the correlations quasi-long ranged. Indeed, as follows from Eq. (10), unless $`g_\theta >g_\theta ^{}2g_\pi /\pi `$, the model is in its “high-temperature” phase with both a spectral gap and exponentially decaying correlations.
One could argue, however, that an ideal sliding phase is “pathological” in the sense of having different time slices completely uncorrelated. In what follows we argue that considering all terms in the classical Hamiltonian (1) does not alleviate the above problem of long-range interactions while bringing new ones.
### 3.2 Non-ideal sliding
One apparent problem arises immediately: according to , in order to have a sliding phase we must have additional gradient couplings between at least next-nearest layers. Since no time derivatives of momenta are allowed to appear in a Hamiltonian, such coupling seems to have no quantum analog. This problem is, however, easily circumvented by doubling the number of components of the field $`\theta `$:
$$𝜽(𝐫,\tau _m)=(\theta ^{(1)}(𝐫,\tau _m),\theta ^{(2)}(𝐫,\tau _m))(\theta _{2m}(𝐫),\theta _{2m+1}(𝐫)).$$
(11)
A single time slice is now represented by two layers. Due to the nearest-layer interactions in the original classical Hamiltonian we have now generated interactions between the two components of our quantum field $`𝜽`$, but this situation is not unusual: a non-linear $`\sigma `$-model provides a standard example of such behaviour. It must now become clear why we have chosen $`p=2`$ in Eq. (4): the Josephson terms now couple only identical components in the nearest time slices.
Keeping track of both components, however, brings unnecessary complications. The problem with constructing a local quantum Hamiltonian is apparent even if we forget about the two-component nature of the field and concentrate only on a single component. In what follows, $`\vartheta `$ will be used to represent either of the two components, $`\theta ^{(1)}`$ or $`\theta ^{(2)}`$, and for simplicity we will only consider the terms in the classical Hamiltonian that do not mix them :
$`\stackrel{~}{H}={\displaystyle \frac{1}{2}}{\displaystyle \underset{n}{}}{\displaystyle }d^2r(K\left[\mathbf{}_{}\theta _n(𝐫)\right]^2+U_2\left\{\mathbf{}_{}[\theta _{n+2}(𝐫)\theta _n(𝐫)]\right\}^2`$
$`2V_\mathrm{J}\mathrm{cos}[\theta _{n+2}(𝐫)\theta _n(𝐫)]).`$ (12)
Denoting $`\vartheta =\theta _n=\theta ^{(i)}(\tau )`$ and $`\vartheta ^{}=\theta _{n+2}=\theta ^{(i)}(\tau +ϵ)`$ we have, similarly to Eq. (6):
$`\vartheta ^{}|\widehat{\stackrel{~}{T}}|\vartheta =\mathrm{exp}[{\displaystyle \frac{\beta }{4}}{\displaystyle }d^2r(K[\mathbf{}\vartheta (𝐫)]^2+K[\mathbf{}\vartheta ^{}(𝐫)]^2+2U_2\{\mathbf{}[\vartheta ^{}(𝐫)\vartheta (𝐫)]\}^2`$
$`4V_\mathrm{J}\mathrm{cos}[\vartheta ^{}(𝐫)\vartheta (𝐫)])].`$ (13)
With the help of $`\mathrm{exp}(a\mathrm{cos}\varphi )=_mI_{|m|}(a)\mathrm{exp}(\mathrm{i}ma)`$, where $`I_n(x)`$ is a modified Bessel function, we obtain, similarly to Eq. (7):
$`\widehat{\stackrel{~}{T}}\mathrm{exp}\left\{{\displaystyle \frac{\beta K}{4}}{\displaystyle d^2r\left[\mathbf{}\widehat{\theta }(𝐫)\right]^2}\right\}`$
$`\times {\displaystyle \underset{\{m(𝐫)\}}{}}\mathrm{exp}\{{\displaystyle }d^2rd^2r^{}\left({\displaystyle \frac{1}{4\pi \beta U_2}}[\widehat{\pi }(𝐫)m(𝐫)][\widehat{\pi }(𝐫^{})m(𝐫^{})]\mathrm{ln}\right|𝐫𝐫^{}|`$
$`+\mathrm{ln}I_{|m(𝐫)|}(\beta V_\mathrm{J}))\}`$
$`\times \mathrm{exp}\left\{{\displaystyle \frac{\beta K}{4}}{\displaystyle d^2r\left[\mathbf{}\widehat{\theta }(𝐫)\right]^2}\right\}.`$ (14)
While obtaining the actual quantum Hamiltonian from Eq. (3.2) is much more complicated than from Eq. (7), the two cases share the same important feature: a long-ranged logarithmic interaction between momenta at different locations.
## 4 Discussion
Let us now discuss the applicability of Hastings’ theorem to the quantum Hamiltonian we have attempted to construct. In short, the theorem states that for a quantum system with a *local* Hamiltonian and a unique ground state separated from excited states by a gap, all equal-time connected correlation functions of local operators $`0|\widehat{A}\widehat{B}|00|\widehat{A}|00|\widehat{B}|0`$ decay exponentially with distance.
Therefore this theorem is not applicable to the quantum Hamiltonians (7,3.2) due to logarithmic interactions appearing there. What seems counterintuitive is that these non-local interactions originate from a perfectly local “sliding” term in the classical action.
The resolution of the original paradox appears to be one aspect of, perhaps, a broader question: when can one go from a local classical Lagrangian to a local quantum Hamiltonian. As our example shows, this needs not always be the case.
We finally remark that locality can sometimes be restored at a cost of introducing auxiliary degrees of freedom. *E.g.*, a “Coulomb gas” analogy for the ideal sliding Hamiltonian (9) readily hints at such a possibility: an introduction of an auxiliary gauge field – an electromagnetic vector potential – will make the theory local. However, this will not contradict Hastings’ theorem, as such field will come with its own gapless mode – a photon!
The authors are grateful to S. Kivelson, S. Chakravarty, A. Kitaev, I. Dimov, X.-G. Wen, E. Fradkin, G. Refael and I. Klich for useful discussions and comments. K. S. and C. N. have been supported by the ARO under Grant No. W911NF-04-1-0236. C. N. has also been supported by the NSF under Grant No. DMR-0411800.
## References |
warning/0506/math0506312.html | ar5iv | text | # Polar actions on symmetric spaces
## 1. Introduction and main results
An isometric action of a compact Lie group on a Riemannian manifold is called polar if there exists a connected immersed submanifold $`\mathrm{\Sigma }`$ which intersects the orbits orthogonally and meets every orbit. Such a submanifold $`\mathrm{\Sigma }`$ is then called a section of the group action. If the section is flat in the induced metric, the action is called hyperpolar. Our main result is a classification of polar actions on compact symmetric spaces with simple isometry group and rank greater that one. This classification shows that these actions are in fact all hyperpolar.
One may think of the elements in a section as being canonical forms, representing the orbits of the group action uniquely up to the action of a finite group, the Weyl group. This point of view may be illustrated by the example of the orthogonal group $`\mathrm{O}(n)`$ acting on the space of real symmetric $`n\times n`$-matrices by conjugation, where the subspace of diagonal matrices is a section. Another motivation comes from submanifold geometry, in particular from the theory of isoparametric submanifolds and their generalizations , . The orbits of polar actions have many remarkable geometric properties, for instance, the principal orbits of polar representations are isoparametric submanifolds of Euclidean space.
However, the history of the subject probably starts with an application in topology. Bott and Bott and Samelson considered the adjoint action of a compact Lie group on itself and on its Lie algebra and more generally, the isotropy action of a compact symmetric space . The motivation of Bott and Samelson to consider these actions was that they are “variationally complete”, which made it possible to apply Morse theory to the space of loops in the symmetric space. Conlon proved that hyperpolar actions on Riemannian manifolds are variationally complete, referring to the sections as $`K`$-transversal domains. Hermann found another class of examples, namely if $`H`$ and $`K`$ are both symmetric subgroups of a simple compact Lie group $`G`$, then the action of $`H`$ on the symmetric space $`G/K`$ is hyperpolar. It was shown much later that actions on compact symmetric spaces are variationally complete if and only if they are hyperpolar.
Conlon observed that s-representations are hyperpolar and later on Dadok obtained a classification of irreducible polar representations. The classification shows that the connected components of the orbits of a polar representation agree with the orbits of an s-representation after a suitable identification of the representation spaces. Reducible polar representations were classified by Bergmann .
Cohomogeneity one actions, i.e. actions whose principal orbits are hypersurfaces, are a special case of independent interest. Cohomogeneity one actions on spheres were classified by Hsiang and Lawson . Later Takagi , D’Atri , and Iwata classified cohomogeneity one actions on $`\mathrm{P}^n`$, $`\mathrm{P}^n`$ and $`𝕆\mathrm{P}^2`$, respectively.
Szenthe , Palais and Terng investigated fundamental properties of polar actions on Riemannian manifolds. Heintze, Palais, Terng and Thorbergsson , obtained structural results for hyperpolar actions on compact symmetric spaces, studied relations to polar actions on infinite dimensional Hilbert space and involutions of affine Kac-Moody algebras. They showed in particular that compact Riemannian homogeneous spaces admitting a hyperpolar action with a fixed point are symmetric.
In , the author gave a classification of hyperpolar actions on the irreducible compact symmetric spaces, the main result being that these actions are orbit equivalent to the examples found by Hermann if the cohomogeneity is $`2`$.
Podestà and Thorbergsson classified polar actions on the compact symmetric spaces of rank one. The first result on polar actions on irreducible symmetric spaces of higher rank without assuming flatness of the sections was obtained by Brück , who showed that polar actions with a fixed point on these spaces are hyperpolar. Podestà and Thorbergsson proved that polar actions on compact irreducible homogeneous Kähler manifolds are coisotropic and classified coisotropic and polar actions on the real Grassmannians $`𝔾_2(^n)`$ of rank two. It turned out that all polar actions on these spaces are hyperpolar.
This approach was further pursued by Biliotti and Gori , who classified coisotropic and polar actions on the complex Grassmannians $`𝔾_k(^n)`$. The classification of coisotropic actions on the compact irreducible Hermitian symmetric spaces was recently completed by Biliotti , showing in particular that polar actions on these spaces are hyperpolar, which led Biliotti to conjecture that this holds for all compact irreducible symmetric spaces.
The present work extends the classification of polar actions to all irreducible symmetric spaces of type I, i.e. to the compact symmetric spaces with simple isometry group, confirming the conjecture of Biliotti for these spaces. We show that polar actions on the symmetric spaces of type I and higher rank are hyperpolar. That is, they are of cohomogeneity one or orbit equivalent to the examples found by Hermann. Our main result can be stated as follows.
###### Theorem 1.
Let $`M`$ be a compact symmetric space of rank greater than one whose isometry group $`G`$ is simple. Let $`HG`$ be a closed connected non-trivial subgroup acting polarly on $`M`$. Then the action of $`H`$ on $`M`$ is hyperpolar, that is, the sections are flat in the induced metric. Moreover, the sections are embedded submanifolds.
In , the hyperpolar actions on irreducible compact symmetric spaces were only determined up to orbit equivalence. In the present work we obtain the complete classification of connected Lie groups acting polarly without fixed points on the symmetric spaces of higher rank with simple isometry group. For actions with fixed points the complete classification follows immediately from Corollary 6.2 and Lemma 2.6.
###### Theorem 2.
Let $`M=G/K`$ be a connected compact symmetric space of rank greater than one whose isometry group is simple. Let $`HG`$ be a closed connected proper subgroup such that the $`H`$-action on $`G/K`$ is polar, non-trivial, non-transitive, and without fixed point. Then
1. either $`HG`$ is maximal connected (and as described in Theorem A of )
2. or the universal cover of the symmetric space $`\stackrel{~}{M}`$ and the conjugacy class of the subgroup $`HG`$ are as given by Table 1, where $`\pi :\mathrm{Isom}(\stackrel{~}{M})G`$ is the covering map, and there exists a connected subgroup $`H_0G`$ whose Lie algebra $`\text{h}_0\text{g}`$ is the fixed point set of an involution of g and such that the $`H_0`$-action on $`G/K`$ has the same orbits as the $`H`$-action.
The first column of Table 1 indicates a connected subgroup $`H_0`$ of $`G`$ containing $`H`$, see Table 3 and the remarks there. By $`\alpha `$ we denote a non-trivial outer automorphism of $`\mathrm{SO}(2n)`$ of order two given by conjugation with an element from $`\mathrm{O}(2n)\mathrm{SO}(2n)`$. The proofs of Theorems 1 and 2 are completed and summarized on pages 1212.
Combining Theorem 1 with the results of and Corollary D of , we obtain the following result on sections and Weyl groups of polar actions. The Weyl group $`W_\mathrm{\Sigma }=\mathrm{N}_H(\mathrm{\Sigma })/Z_H(\mathrm{\Sigma })`$ is a quotient group of the group $`\widehat{W}_\mathrm{\Sigma }`$ as defined in Lemma 5.1.
###### Corollary 1.
Let $`H`$ be a connected compact Lie group acting polarly on a compact symmetric space $`M`$ with simple isometry group. Then a section $`\mathrm{\Sigma }`$ of the $`H`$-action on $`M`$ is isometric to a flat torus, a sphere or a real projective space. The group $`\widehat{W}_\mathrm{\Sigma }`$ acting on the universal cover of $`\mathrm{\Sigma }`$ is an irreducible affine Coxeter group in case $`\mathrm{\Sigma }`$ is flat or a finite Coxeter group of Euclidean space restricted to a sphere in case $`\mathrm{\Sigma }`$ is non-flat.
In particular, the Weyl groups of such polar actions can be described by connected Dynkin diagrams of affine type (in the hyperpolar case) or Dynkin diagrams of the finite type (in the polar, non-hyperpolar case).
This article is organized as follows. We start by setting up terminology and notation. We then review examples and known results on polar actions. In Section 3 we recall some facts about symmetric spaces and their totally geodesic submanifolds; in particular, we give a characterization of maximal totally geodesic submanifolds and obtain an upper bound on the dimension of totally geodesic submanifolds locally isometric to a product of spheres.
In Section 4, we recall a criterion which reduces the problem of deciding whether an action on a symmetric spaces is polar or not to a problem on the Lie algebra level. In Section 5 we prove the Splitting Theorem 5.2 which says that if a section $`\mathrm{\Sigma }`$ of a polar action admits a local splitting $`\stackrel{~}{\mathrm{\Sigma }}=\stackrel{~}{\mathrm{\Sigma }}_1\times \stackrel{~}{\mathrm{\Sigma }}_2`$ such that the Weyl group acts trivially on one factor $`\stackrel{~}{\mathrm{\Sigma }}_2`$, then the symmetric space is locally a Riemannian product $`M\times \stackrel{~}{\mathrm{\Sigma }}_2`$. As a consequence, we show that the section of a polar action on a compact irreducible symmetric space is locally isometric to a product of spaces of constant curvature. This observation is crucial for our classification since it implies an upper bound on the cohomogeneity, reducing the classification problem to a finite number of cases.
In Section 6 we introduce another main tool by collecting various sufficient conditions for actions to be polarity minimal, which means that the restriction to a closed connected subgroup with orbits of lower dimension is either non-polar or trivial. This is of essential importance since it enables us to restrict our attention at first to maximal subgroups of the isometry group. In many cases we are able to show that the action of a maximal connected subgroup is non-polar and polarity minimal, thereby excluding all of its subgroups.
In the remaining part of the paper, the classification is carried out. We start with the maximal connected subgroups in the isometry group of a symmetric space. In Section 7, we consider Hermann actions, i.e. actions of symmetric subgroups of the isometry group. We show that actions of cohomogeneity $`2`$ are polarity minimal and determine orbit equivalent subactions. We then consider maximal connected subgroups in the isometry group of classical symmetric spaces which are given by irreducible representations of non-simple groups. It turns out that they are either non-polar and polarity minimal or of cohomogeneity one. In Section 9 we study actions of simple irreducible subgroups in the classical groups. In Section 10 we consider actions on the exceptional symmetric spaces. It turns out that the actions of non-symmetric maximal subgroups are non-polar and polarity minimal.
It then remains to study subactions of cohomogeneity one and transitive actions. Since we do not have an a priori proof that these actions are polarity minimal, it is necessary to descend from maximal connected subgroups $`H_1G`$ acting with cohomogeneity $`1`$ to further subgroups
$$GH_1H_2\mathrm{},$$
where $`H_{n+1}H_n`$ is maximal connected, until we arrive at an action which is polarity minimal.
#### Acknowledgements
I would like to thank Ernst Heintze, Mamoru Mimura, Chuu-Lian Terng and Gudlaugur Thorbergsson for helpful discussions and comments; I am especially indebted to Burkhard Wilking for providing a crucial step in the proof of the Splitting Theorem 5.2 and for pointing out an error in an earlier version of this paper.
## 2. Preliminaries and examples
An isometric action of a compact Lie group $`G`$ on a Riemannian manifold is called polar if there exists a connected immersed submanifold $`\mathrm{\Sigma }`$ such that $`\mathrm{\Sigma }`$ meets all $`G`$-orbits and the intersection of $`\mathrm{\Sigma }`$ with any $`G`$-orbit is orthogonal at all intersection points. Such a submanifold $`\mathrm{\Sigma }`$ is called a section for the $`G`$-action on $`M`$. In particular, the actions of finite groups and transitive actions are special cases of polar actions, the section being the whole space or a point, respectively.
Note that we do not require the section to be an embedded submanifold, generalizing the definitions of and . However, it turns out by our classification that on symmetric spaces of type I the sections are flat and therefore closed embedded submanifolds by Corollary 2.12 of .
The dimension of $`\mathrm{\Sigma }`$ equals the cohomogeneity of the $`G`$-action and hence the tangent space $`\mathrm{T}_p\mathrm{\Sigma }`$ at a regular point $`p\mathrm{\Sigma }`$ coincides with the normal space $`\mathrm{N}_p(Gp)`$ at $`p`$ of the orbit through $`p`$. From this, it follows that any two sections are mapped isometrically onto each other by some group element. It has been proved in that sections are totally geodesic.
In the special case where the sections are flat in the induced metric, the action is called hyperpolar. Examples for hyperpolar actions are given by the action of a compact Lie group on itself by conjugation, where the sections are the maximal tori. More generally, the action of an isotropy group of a symmetric space is hyperpolar, the sections being the flats of the symmetric space.
For a polar action one can define the Weyl group by considering the normalizer of a section, i.e. all group elements which map the section onto itself, this group acts on the section by isometries and the Weyl group is defined by factoring out the kernel of this action.
###### Definition 2.1.
Let $`M`$ be a Riemannian manifold on which the compact Lie group $`G`$ acts polarly with section $`\mathrm{\Sigma }`$. The (generalized) Weyl group $`W_\mathrm{\Sigma }=W_\mathrm{\Sigma }(M,G)`$ is the group $`N_G(\mathrm{\Sigma })/Z_G(\mathrm{\Sigma })`$, where $`N_G(\mathrm{\Sigma })=\{gGg\mathrm{\Sigma }=\mathrm{\Sigma }\}`$ and $`Z_G(\mathrm{\Sigma })=\{gGgs=s\text{ for all}s\mathrm{\Sigma }\}`$ are the normalizer and centralizer of $`\mathrm{\Sigma }`$ in $`G`$, respectively.
Two Riemannian $`G`$-manifolds are called conjugate if there exists an equivariant isometry between them. In particular, the actions of two conjugate subgroups of the isometry group of a Riemannian manifold are conjugate. To study isometric actions on a Riemannian manifold it suffices to consider conjugacy classes of subgroups in the isometry group.
Two isometric actions of two Lie groups $`G`$ and $`G^{}`$ on a Riemannian manifold $`M`$ are called orbit equivalent if there exists an isometry of $`M`$ which maps $`G`$-orbits onto $`G^{}`$-orbits; they are called locally orbit equivalent if there is an isometry mapping connected components of $`G`$-orbits onto connected components of $`G^{}`$-orbits. Obvious examples of orbit equivalent actions are given by various groups acting transitively on spheres, e.g. the actions of $`\mathrm{SO}(4n)`$, $`\mathrm{SU}(2n)`$, $`\mathrm{U}(2n)`$, and $`\mathrm{Sp}(n)`$ on $`^{4n}`$ are all orbit equivalent.
We use the term subaction for the restriction of an action of a group $`G`$ to a subgroup $`HG`$; in case the $`H`$-orbits coincide with the $`G`$-orbits, the $`H`$-action is called orbit equivalent subaction.
A normal subgroup $`N`$ of a compact Lie group $`G=G^{}N`$ acting isometrically on a Riemannian manifold is called inessential if the $`G`$-action restricted to $`G^{}`$ is orbit equivalent to the $`G`$-action. An isometric action of a compact connected Lie group $`G`$ on a Riemannian manifold $`M`$ is called orbit maximal if any other isometric action of any other compact connected Lie group $`G^{}`$ such that every $`G`$-orbit is contained in a $`G^{}`$-orbit is either orbit equivalent or transitive on $`M`$.
An immersed submanifold $`M`$ in a symmetric space $`N`$ is said to have parallel focal structure if the normal bundle $`\nu (M)`$ is globally flat and the focal data is invariant under normal parallel translation, that is, for every parallel normal field $`v`$ on $`M`$ the rank of $`d\eta _{v(x)}`$ is locally constant on $`M`$, where the end point map $`\eta :\mathrm{N}MN,v\mathrm{exp}(v)`$ is defined to be the restriction of the exponential map to the normal bundle $`\nu (M)`$, see . The principal orbits of a polar action on a symmetric space have parallel focal structure .
A submanifold with parallel focal structure is called equifocal if the normal bundle $`\nu (M)`$ is abelian, that is, $`\mathrm{exp}(\nu (M))`$ is contained in some totally geodesic flat subspace of $`N`$ for each point $`xM`$. The principal orbits of hyperpolar actions on symmetric spaces of compact type are equifocal submanifolds, see , Theorem 2.1. Our Theorem 1 shows that submanifolds with parallel focal structure which arise as principal orbits of polar actions on symmetric spaces of higher rank with simple compact isometry group are in fact equifocal. We conjecture that more generally submanifolds with parallel focal structure in irreducible compact symmetric spaces of higher rank are equifocal, hence of codimension one or homogeneous by the result of Christ .
### 2.1. Notation
We will frequently use the following notational conventions for compact Lie groups and their representations. We view the classical Lie groups $`\mathrm{SO}(n)`$, $`\mathrm{SU}(n)`$, and $`\mathrm{Sp}(n)`$ as matrix Lie groups as described in , Ch. X, § 2.1. We assume that reducible subgroups of the classical groups are standardly embedded, e.g. by $`\mathrm{SO}(m)\times \mathrm{SO}(n)`$ we denote the subgroup
(2.1)
$$\left\{\left(\begin{array}{cc}A& \\ & \\ & B\end{array}\right)\right|A\mathrm{SO}(m),B\mathrm{SO}(n)\}\mathrm{SO}(m+n).$$
We write $`H_1H_2`$ for the Kronecker product of two matrix Lie groups. When we write $`\mathrm{G}_2`$, we refer to an irreducible representation by orthogonal $`7\times 7`$-matrices; similarly, $`\mathrm{Spin}(7)`$ stands for a matrix Lie group which is the image of the $`8`$-dimensional spin representation of $`\mathrm{Spin}(7)`$. By $`^n`$, $`^n`$, $`^n`$ we will denote the standard representation of $`\mathrm{O}(n)`$, $`\mathrm{U}(n)`$ or $`\mathrm{Sp}(n)`$, respectively.
### 2.2. Polar representations
An important class of examples for polar actions is given by polar representations on Euclidean space. Since the sections of polar actions are totally geodesic, they are linear subspaces in the case of polar representations and polar representations are therefore automatically hyperpolar. Polar representations are of importance for our classification since they occur as slice representations of polar actions. Let $`M`$ be a Riemannian $`G`$-manifold and let $`G_p`$ be the isotropy subgroup at $`p`$. The restriction of the isotropy representation to $`\mathrm{N}_p(Gp)`$ is called the slice representation at $`p`$. Slice representations are a fundamental tool for the study of Lie group actions since they provide a means to describe the local behavior of an action in a tubular neighborhood of an orbit by a linear representation.
###### Slice Theorem 2.2.
Let $`M`$ be a Riemannian $`G`$-manifold, let $`pM`$ and $`V=\mathrm{N}_p(Gp)`$ the normal space at $`p`$ to the $`G`$-orbit through $`p`$. Then there is an equivariant diffeomorphism $`\mathrm{\Psi }`$ of a $`G`$-invariant open neighborhood around the zero section in the normal bundle $`G\times _{G_p}VG/G_p`$ onto a $`G`$-invariant open neighborhood around the orbit $`Gp`$ such that the zero section in $`G\times _{G_p}V`$ is mapped to the orbit $`Gp`$. The diffeomorphism $`\mathrm{\Psi }`$ is given by the end point map which maps any normal vector $`v_qN_q(Gp)`$ to its image under the exponential map $`\mathrm{exp}_q(v_q)`$.
###### Proof.
See e.g. , p. 3. ∎
It is an immediate consequence of the Slice Theorem 2.2 that the slice representation and the $`G`$-action on $`M`$ have the same cohomogeneity. Slice representations are in particular useful for our classification since the polarity of an action is inherited by its slice representations.
###### Proposition 2.3.
Let $`M`$ be a Riemannian $`G`$-manifold. If the action on $`N`$ is polar then for all $`pM`$ the slice representation at $`p`$ is polar with $`\mathrm{T}_p\mathrm{\Sigma }`$ as a section, where $`\mathrm{\Sigma }`$ is the section of the $`G`$-action on $`M`$ containing $`p`$.
###### Proof.
This was proved in , Theorem 4.6. Although in the sections are assumed to be embedded submanifolds, the proof is still valid if one requires the sections only to be immersed. ∎
We use the term effectivized slice representation to describe the representation of the isotropy group with the effectivity kernel factored out. Let us recall some known results about polar representations.
###### Definition 2.4.
Let $`G`$ be a compact Lie group and $`K`$ be a closed subgroup. By
$$\chi (G,K)=\mathrm{Ad}_G|_K\mathrm{Ad}_K$$
we denote the equivalence class of the isotropy representation of the homogeneous space $`G/K`$, i.e. the restriction of the adjoint representation of $`G`$ to $`K`$ acting on a $`K`$-invariant complement of k in g. In the special case of a symmetric pair $`(G,K)`$, see below, the (equivalence class of the) representation $`\chi (G/K)`$ is called an s-representation.
A compact subgroup $`K`$ of a Lie group $`G`$ is called symmetric subgroup if there exists an involutive automorphism of $`G`$ such that $`G_0^\sigma KG^\sigma `$, where $`G^\sigma `$ and $`G_0^\sigma `$ denote the fixed point set of $`\sigma `$ and its connected component, respectively. A pair $`(G,K)`$, where $`G`$ is a Lie group and $`K`$ a symmetric subgroup is called a symmetric pair. Any Riemannian globally symmetric space $`M`$ has a homogeneous presentation $`G/K`$, where $`G`$ is the isometry group of $`M`$, such that $`K`$ is a symmetric subgroup of $`G`$. Conversely, if $`(G,K)`$ is a symmetric pair, then $`G/K`$ endowed with a $`G`$-invariant metric is a Riemannian globally symmetric space, see .
It is well known that the adjoint representations of compact Lie groups and more generally s-representations are polar. As far as concerns the geometry of the orbits, also the converse is true.
###### Theorem 2.5 (Dadok).
A representation $`\rho :G\mathrm{O}(n)`$ of a compact Lie group is polar if and only if it is locally orbit equivalent to an s-representation, i.e. the connected components of its orbit agree with the orbits of an s-representation after a suitable isometric identification of the representation spaces.
###### Proof.
The proof given in relies on a classification of the irreducible polar representations. See for a conceptual proof in case the cohomogeneity is $`3`$. See for an alternative proof, where a similar classification strategy as in the present work is used. ∎
It is shown in Theorem 3.12 of that irreducible polar representations of cohomogeneity $`2`$ are orbit maximal when restricted to a sphere around the origin. For irreducible s-representations of cohomogeneity $`2`$, orbit equivalent subgroups were determined in . We state the result below.
###### Lemma 2.6.
Let $`G`$ be a connected simple compact Lie group and let $`K`$ be a connected symmetric subgroup such that $`\mathrm{rk}(G/K)2`$. Let $`HK`$ be a closed connected subgroup. Then $`\chi (G,K)`$ and $`\chi (G,K)|_H`$ are orbit equivalent if and only if either $`H=K`$ or the triple $`(G,K,H)`$ is as given in Table 2.
### 2.3. Hyperpolar actions on symmetric spaces
If $`H`$, $`K`$ are two symmetric subgroups of the compact Lie group $`G`$, then the action of $`H`$ on $`G/K`$ is hyperpolar . Slightly more generally, if $`H`$ is a subgroup of $`G`$ such that its Lie algebra $`\text{h}\text{g}`$ is the fixed point set of an involution of g, then the action of $`H`$ on the symmetric space $`G/K`$ is hyperpolar and we call such actions Hermann actions. In the special case $`H=K`$ we have the isotropy action of the symmetric space and the sections are just the flats of the symmetric space. It was shown in that all hyperpolar actions on irreducible symmetric spaces of compact type are of cohomogeneity one or orbit equivalent to Hermann actions.
All fixed-point free Hermann actions on the symmetric spaces of type I are given by Table 3. Here $`\alpha `$ denotes the non-trivial diagram automorphism of $`\mathrm{SO}(2n)`$ given by conjugation with a matrix from $`\mathrm{O}(2n)\mathrm{SO}(2n)`$ and $`\tau `$ an order three diagram automorphism of $`\mathrm{Spin}(8)`$. The type of the Hermann action indicated in the first column refers to the type of the symmetric subgroups involved as given in Table 4, e.g. the symbol A I-II refers to the action of $`H`$ on $`G/K`$, where $`G/H=\mathrm{SU}(2n)/\mathrm{SO}(2n)`$ is a symmetric space of type A I and $`G/K=\mathrm{SU}(2n)/\mathrm{Sp}(n)`$ is a symmetric space of type A II; whereas for the action $`K`$ on $`G/H`$ we use the notation A II-I. For the conjugacy classes of connected symmetric subgroups in simple compact Lie groups see , 3.1.1 and 3.1.2. The cohomogeneity of the actions is given in the last column.
Hyperpolar actions on compact symmetric spaces have the remarkable property that they lift under certain Riemannian submersions to actions which are again hyperpolar, cf. .
###### Proposition 2.7.
Let $`G`$ be a compact simple Lie group and let $`KG`$ be a symmetric subgroup, $`M=G/K`$ the corresponding symmetric space and $`H`$ a closed subgroup of $`G`$. Then the $`H`$-action on $`M`$ is hyperpolar if and only if the $`H\times K`$-action on $`G`$ is hyperpolar.
###### Proof.
See , Proposition 2.11 ∎
It can be shown using Proposition 4.1 that polar actions on compact symmetric spaces have this lifting property only if they are hyperpolar. In particular, if we lift the known polar actions on symmetric spaces to the groups, we do not obtain any examples of polar actions besides the hyperpolar ones. Another remarkable property of hyperpolar actions is that they are orbit maximal on irreducible symmetric spaces of compact type.
###### Proposition 2.8.
Let $`M=G/K`$ be a connected irreducible symmetric space of compact type and $`HLG`$ closed connected subgroups. If the $`H`$-action on $`M`$ is hyperpolar, then the $`L`$-action on $`M`$ is transitive or orbit equivalent to the $`H`$-action.
###### Proof.
See , Corollary D. ∎
The non-orbit maximal examples of polar actions found by Podestà and Thorbergsson show that Proposition 2.8 does not directly generalize to polar actions, see below. However, it is a consequence of our classification that polar actions on symmetric spaces of rank $`2`$ with simple compact isometry group are orbit maximal.
### 2.4. Polar actions on rank one symmetric spaces
Polar actions on rank one symmetric spaces have been classified by Podestà and Thorbergsson . The hyperpolar, i.e. cohomogeneity one, actions on these spaces had before been classified in , , and . The results can be summarized as follows. The classification of polar actions on spheres (and real projective spaces) follows from , since every polar action on the sphere is given as the restriction of a polar representation to the sphere. The isotropy representations of Hermitian symmetric spaces of real dimension $`2n+2`$ induce polar actions on $`\mathrm{P}^n`$ and all polar actions on $`\mathrm{P}^n`$ are orbit equivalent to actions obtained in this fashion. Similarly, all polar actions on $`\mathrm{P}^n`$ come from isotropy representations of products of quaternion-Kähler symmetric spaces, with the additional restriction that all factors but one must be of rank one. While all these polar actions arise from polar actions on the sphere, the actions on the Cayley plane $`𝕆\mathrm{P}^2=\mathrm{F}_4/\mathrm{Spin}(9)`$ do not have such an interpretation. The maximal connected subgroup $`\mathrm{SU}(3)\mathrm{SU}(3)\mathrm{F}_4`$ acts polarly on the Cayley plane with cohomogeneity two. The groups $`\mathrm{Sp}(3)\mathrm{Sp}(1)`$, $`\mathrm{Sp}(3)\mathrm{U}(1)`$, $`\mathrm{Sp}(3)`$ and $`\mathrm{Spin}(9)`$ act with cohomogeneity one. In addition there are three polar actions of cohomogeneity two with a fixed point of the following subgroups of $`\mathrm{Spin}(9)`$:
$$\mathrm{Spin}(8),\mathrm{SO}(2)\mathrm{Spin}(7),\mathrm{Spin}(3)\mathrm{Spin}(6).$$
In particular, polar actions on rank one symmetric spaces are not orbit maximal in general.
## 3. Symmetric spaces and their totally geodesic submanifolds
In the following we will collect some useful facts about symmetric spaces and their totally geodesic submanifolds. Sections of polar actions are totally geodesic submanifolds and it will be shown in Theorem 5.4 that the sections of a non-trivial polar action on an irreducible compact symmetric space are locally isometric to Riemannian products whose factors are spaces of constant curvature. We give an upper bound on the dimension of such submanifolds in Lemma 3.3. For the proof of Theorem 5.4, which is essentially a consequence of the Splitting Theorem 5.2, we will need the characterization of totally geodesic hypersurfaces in reducible symmetric spaces given in Corollary 3.5, because the Weyl group of a polar action is generated by reflections in totally geodesic hypersurfaces. We will conclude this section by recalling a well known characterization of maximal subgroups in the classical groups.
Every symmetric space $`M`$ may be presented as $`G/K`$, where $`G`$ is the isometry group of $`M`$ and $`K`$ is a symmetric subgroup of $`G`$. Conversely, if $`(G,K)`$ is a symmetric pair, then $`G/K`$ is a symmetric space if it is equipped with an appropriate metric. A Riemannian manifold is an irreducible Riemannian symmetric space of compact type if and only if it is isometric to
* either $`G/K`$, where $`G`$ is a simple, compact, connected Lie group and $`K`$ a symmetric subgroup of G (symmetric space of type I)
* or a simple, compact, connected Lie group equipped with a biinvariant metric (symmetric space of type II)
The local isometry classes of the symmetric spaces of type I are given by Table 4. By $`\mathrm{SO}^{}(2n)`$ we denote the image of a half-spin representation of $`\mathrm{Spin}(2n)`$.
The global isometry classes of symmetric spaces are given by the following theorem, which will be needed for the proof of the Splitting Theorem 5.2.
###### Theorem 3.1.
Let $`M`$ be a simply connected Riemannian symmetric space with decomposition $`M=M_0\times M_1\times \mathrm{}\times M_t`$ into Euclidean and irreducible parts. Define $`G=V\times I(M_1)_0\times \mathrm{}\times I(M_t)_0`$ where $`V`$ is the vector group of pure translations of the Euclidean space $`M_0`$. Define $`\mathrm{\Delta }=V\times \mathrm{\Delta }_1\times \mathrm{}\times \mathrm{\Delta }_t`$ where $`\mathrm{\Delta }_i`$ ($`i>0`$) is the centralizer of $`I(M_i)_0`$ in $`I(M_i)`$. Then $`G`$ is the group generated by all transvections of $`M`$ and $`\mathrm{\Delta }`$ is the centralizer of $`G`$ in $`I(M)`$. In particular, the symmetric spaces covered by $`M`$ are just the manifolds $`M/\mathrm{\Gamma }`$ where $`\mathrm{\Gamma }`$ is a discrete subgroup of $`\mathrm{\Delta }`$.
The group $`\mathrm{\Delta }_i`$ is trivial if $`M_i`$ is noncompact, and is finite if $`M_i`$ is compact. In particular, the discrete subgroups of $`\mathrm{\Delta }`$ are just the subgroups $`\mathrm{\Gamma }\mathrm{\Delta }`$ with discrete projection on the vector group $`V`$.
###### Proof.
See , Ch. 8, Sec. 3. ∎
Totally geodesic submanifolds of symmetric spaces correspond to Lie triple systems.
###### Proposition 3.2.
Let $`M`$ be a Riemannian globally symmetric space and let $`p_0M`$. Let $`G=I(M)_0`$ and let $`K=G_{p_0}`$. Let $`\text{g}=\text{k}\text{p}`$, where k is the Lie algebra of $`K`$ and where we identify $`\text{p}=\mathrm{T}_{p_0}M`$ as usual. Let $`\sigma _{}:\text{g}\text{g}`$ be the automorphism of g which acts on k as $`\mathrm{id}_\text{k}`$ and on p as $`\mathrm{id}_\text{p}`$. The totally geodesic submanifolds of $`M`$ containing $`p_0`$ are in one-to-one correspondence with the Lie triple systems $`𝔰\text{p}`$; i.e. if $`𝔰\text{p}`$ is a Lie triple system, then $`\mathrm{exp}(𝔰)M`$ is a totally geodesic submanifold and, conversely, if $`SM`$ is a totally geodesic submanifold, then $`\mathrm{T}_{p_0}\mathrm{\Sigma }\text{p}`$ is a Lie triple system.
Moreover, for any Lie triple system $`𝔰\text{p}`$, define $`\text{g}^{}=𝔰+[𝔰,𝔰]`$ and $`\text{k}^{}=[𝔰,𝔰]`$; then $`\text{g}^{}`$ is the Lie subalgebra of g generated by $`𝔰`$, $`\text{g}^{}`$ is invariant under $`\sigma _{}`$ and $`\text{k}^{}=\text{g}^{}\text{k}`$. Let $`G^{}`$ and $`K^{}`$ be the connected Lie subgroups of $`G`$ with Lie algebras $`\text{g}^{}`$ and $`\text{k}^{}`$, respectively. Then $`(G^{},K^{})`$ is a symmetric pair and $`G^{}`$ acts transitively on $`\mathrm{exp}(𝔰)`$.
###### Proof.
See , Ch. IV, § 7. ∎
###### Lemma 3.3.
Let $`(G,K)`$ be a symmetric pair such that $`M=G/K`$ is a Riemannian symmetric space of compact type and let $`\mathrm{\Sigma }M`$ be a totally geodesic submanifold whose universal cover is a product of spheres. Then $`dim(\mathrm{\Sigma })\mathrm{rk}(G)+\mathrm{rk}(K)`$.
###### Proof.
Let$`(G_\mathrm{\Sigma },K_\mathrm{\Sigma })`$ be the symmetric pair corresponding to $`\mathrm{\Sigma }`$, we have $`\text{g}_\mathrm{\Sigma }=\text{g}_\mathrm{\Sigma }^1\mathrm{}\text{g}_\mathrm{\Sigma }^m`$, where $`\text{g}_\mathrm{\Sigma }^i𝔰𝔬(n_i+1)`$. Let $`\text{k}_\mathrm{\Sigma }^i\text{g}_\mathrm{\Sigma }^i`$ such that $`\text{k}_\mathrm{\Sigma }^i𝔰𝔬(n_i)`$. Let $`\text{g}=\text{k}\text{p}`$ as usual. By Proposition 3.2, we may assume $`\text{k}_\mathrm{\Sigma }^i\text{k}`$. Now choose maximal abelian subalgebras $`𝔞_i\text{g}_\mathrm{\Sigma }^i`$ as follows. If $`n_i`$ is even, then $`\mathrm{rk}(\text{k}_\mathrm{\Sigma }^i)=\mathrm{rk}(\text{g}_\mathrm{\Sigma }^i)`$ and we may choose $`𝔞^i=𝔞_i^\text{k}\text{k}_\mathrm{\Sigma }^i`$. If $`n_i`$ is odd, then we may choose $`𝔞^i\text{g}_\mathrm{\Sigma }^i`$ such that $`𝔞_i=𝔞_i^\text{p}𝔞_i^\text{k}`$ where $`𝔞_i^\text{p}\text{p}`$ is one-dimensional and $`𝔞_i^\text{k}\text{k}_\mathrm{\Sigma }^i`$. Let $`𝔞^\text{k}=_{i=1}^m𝔞_i^\text{k}`$ and $`𝔞^\text{p}=_{n_i1(2)}𝔞_i^\text{p}`$. Then we have $`dim(\mathrm{\Sigma })=2dim(𝔞^\text{k})+dim(𝔞^\text{p})`$. Since $`𝔞^\text{k}𝔞^\text{p}`$ is an abelian subalgebra of g and $`𝔞^\text{k}`$ is an abelian subalgebra of k, it follows that $`dim(\mathrm{\Sigma })\mathrm{rk}(G)+\mathrm{rk}(K)`$. ∎
The estimate on the dimension given by the Lemma above is not optimal in all cases. See for classifications of totally geodesic submanifolds in symmetric spaces.
The following Theorem, which characterizes maximal totally geodesic submanifolds of reducible symmetric spaces, is an analogue of Theorem 15.1 in , which characterizes maximal subalgebras of semisimple Lie algebras; we give a proof which is similar to the proof in .
###### Theorem 3.4.
Let $`S`$ be a connected simply connected symmetric space with decomposition $`S=S_0\times S_1\times \mathrm{}\times S_k`$ such that $`S_1,\mathrm{},S_k`$ are irreducible and $`S_0`$ is of Euclidean type. Let $`V`$ be a maximal totally geodesic submanifold of $`S`$, (i.e. if there is a totally geodesic submanifold $`W`$ such that $`VWS`$ then either $`V=W`$ or $`W=S`$). Let $`p=(p_0,\mathrm{},p_k)V`$. Then either there is an index $`i\{0,\mathrm{},k\}`$ and a totally geodesic submanifold $`\stackrel{~}{V}S_i`$ such that
$$V=S_0\times \mathrm{}\times S_{i1}\times \stackrel{~}{V}\times S_{i1}\times \mathrm{}\times S_k,$$
or there are two factors $`S_i`$ and $`S_j`$ ($`ij`$) and a map $`\varphi :S_iS_j`$ which is an isometry up to scaling such that
$$V=\underset{\genfrac{}{}{0pt}{}{\mathrm{}=1}{\mathrm{}i,j}}{\overset{k}{}}S_{\mathrm{}}\times \{(x,\varphi (x))|xS_i\}.$$
###### Proof.
Let $`G=I(S)`$ and let $`K=I(S)_p`$ such that $`\text{g}=\text{k}\text{p}`$ is a Cartan decomposition associated with the symmetric space $`S=G/K`$. Let $`G_i=I(S_i)`$ and let $`K_i=I(S_i)_{p_i}`$ such that $`\text{g}_i=\text{k}_i\text{p}_i`$ are Cartan decompositions corresponding to the irreducible factors $`S_i=G_i/K_i`$. Since $`V`$ is a totally geodesic submanifold, we have that $`\nu =\mathrm{T}_pV\text{p}`$ is a Lie triple system by Proposition 3.2. Obviously, the projection $`\pi _i(\nu )`$ onto each of the summands $`\text{p}_i`$ is again a Lie triple system. Now there are two cases:
Either there is an index $`i\{0,\mathrm{},k\}`$ such that $`\mathrm{pr}_i(V)S_i`$, where $`\mathrm{pr}_i:SS_i`$ denotes the canonical projection onto $`S_i`$. Then $`\mathrm{pr}_i(V)`$ is a totally geodesic submanifold in $`S_i`$ and there is a maximal totally geodesic submanifold $`\stackrel{~}{V}\text{p}_i`$ containing $`\mathrm{pr}_i(V)`$. Thus, $`S_0\times \mathrm{}\times S_{i1}\times \stackrel{~}{V}\times S_{i1}\times \mathrm{}\times S_k`$ is a totally geodesic submanifold of $`S`$ containing $`V`$ and which is, by maximality, equal to $`V`$.
Or $`\pi _i(\nu )=\text{p}_i`$ for all $`i=0,\mathrm{},k`$, where the Lie algebra epimorphisms $`\pi _i:\text{g}\text{g}_i`$ are given by the canonical projections. In this case, it follows that there are at least two indices $`i,j\{0,\mathrm{},k\}`$ such that $`\text{p}_i`$ and $`\text{p}_j`$ are both not contained in $`\nu `$. Define $`\nu ^{}=\nu (\text{p}_i\text{p}_j)\nu `$. This is a Lie triple system in p, since it is the intersection of two Lie triple systems. Hence $`_{\genfrac{}{}{0pt}{}{\mathrm{}=0}{\mathrm{}i,j}}^k\text{p}_{\mathrm{}}\nu ^{}`$ is a Lie triple system in p which contains $`\nu `$ and is different from p, thus, by maximality, is the tangent space $`\mathrm{T}_pV`$.
It remains now to study the Lie triple system $`\nu ^{}\text{p}_i\text{p}_j`$. By Proposition 3.2, it follows that the Lie algebra $`\text{g}^{}=\nu ^{}[\nu ^{},\nu ^{}]`$ generated by $`\nu ^{}`$ is the Lie algebra of a group $`G^{}G_i\times G_j`$ acting transitively on the totally geodesic submanifold $`V^{}`$ of $`S_i\times S_j`$ which is the exponential image of $`\nu ^{}\mathrm{T}_{(p_i,p_j)}(S_i\times S_j)`$.
We show that $`\text{g}^{}\text{g}_i`$ is an ideal in $`\text{g}_i`$: Let $`x\text{g}_i`$, $`y\text{g}^{}\text{g}_i`$, there is a $`z\text{g}^{}`$ such that $`\pi _i(z)=x`$ and it follows that $`[x,y]=[z,y]\text{g}^{}\text{g}_i`$. By the same argument, $`\text{g}^{}\text{g}_j`$ is an ideal in $`\text{g}_j`$.
Let us assume for the moment that $`i,j0`$. Since $`\pi _i(\nu )=\text{p}_i`$ and $`\pi _j(\nu )=\text{p}_j`$, we have that $`\pi _i(\text{g}^{})=\text{g}_i`$ and $`\pi _j(\text{g}^{})=\text{g}_j`$ since $`S_i`$ and $`S_j`$ are irreducible symmetric spaces. Since they are the Lie algebras of isometry groups of irreducible symmetric spaces, the $`\text{g}_i`$ are either simple or the direct sum of two isomorphic simple ideals $`\text{h}_i\text{h}_i`$ (in case $`S_i`$ is of type II). Therefore the ideal $`\text{g}^{}\text{g}_i`$ is either zero, equal to $`\text{g}_i`$ or equal to $`\text{h}_i`$. The last case is impossible, since $`\text{g}^{}`$ has to be invariant under the action of the Cartan involution of $`\text{h}_i\text{h}_i`$, which is given by $`(x,y)(y,x)`$; the case $`\text{g}^{}\text{g}_i=\text{g}_i`$ is also impossible, since $`\text{p}_i`$ is not contained in $`\nu ^{}`$. Now we can show that $`\text{g}_i`$ and $`\text{g}_j`$ are isomorphic: Let $`x\text{g}_i`$ then there is a an element $`y\text{g}_j`$ such that $`(x,y)\text{g}^{}`$; but this element is uniquely defined since otherwise, $`\text{g}_j`$ would have a non-trivial intersection with $`\text{g}^{}`$. The map $`\text{g}_i\text{g}_j`$ we defined in this way is easily seen to be a Lie algebra isomorphism and the subalgebra $`\text{g}^{}`$ is given by the diagonal embedding of $`\text{g}^{}\text{g}^{}\text{g}^{}\text{g}_i\text{g}_j`$. It remains to be shown that the spaces $`S_i`$ and $`S_j`$ are isometric up to scaling: This follows from the requirement that the Cartan involution corresponding to $`\text{g}_i\text{g}_j=(\text{k}_i\text{k}_j)(\text{p}_i\text{p}_j)`$ must leave the diagonally embedded subalgebra $`\text{g}^{}`$ invariant and is thus of the form $`(x,y)(\sigma (x),\sigma (y))`$, where $`\sigma `$ is an involution of $`\text{g}_i\text{g}_j`$.
Finally, assume $`i=0`$, $`j\{1,\mathrm{},k\}`$. This case can be included in the above proof if we further split $`\text{p}_0`$ into the direct sum of $`\text{g}^{}\text{p}_0`$ plus a complementary subspace. Then one is again in the situation that g can be written as a direct sum of ideals all of which have either trivial intersection with the Lie algebra generated by $`\nu `$ or are contained in this Lie algebra. Then the same type of argument leads to the contradiction that an abelian Lie algebra is isomorphic to one of $`\text{g}_1,\mathrm{},\text{g}_k`$. ∎
###### Corollary 3.5.
Let $`S`$ be a connected simply connected symmetric space with decomposition $`S=S_0\times S_1\times \mathrm{}\times S_k`$ such that $`S_1,\mathrm{},S_k`$ are irreducible and $`S_0`$ is of Euclidean type. Let $`H`$ be a totally geodesic hypersurface of $`S`$. Let $`p=(p_0,\mathrm{},p_k)H`$. Then there is an index $`i\{0,\mathrm{},k\}`$ and a totally geodesic hypersurface $`\stackrel{~}{H}S_i`$ such that
$$H=S_0\times \mathrm{}\times S_{i1}\times \stackrel{~}{H}\times S_{i1}\times \mathrm{}\times S_k.$$
###### Proof.
Obviously, a totally geodesic hypersurface is a maximal totally geodesic submanifold, so we may apply Theorem 3.4. Irreducible non-flat symmetric spaces are at least of dimension two, thus the second possibility in the assertion of Theorem 3.4 does not occur here, since this would lead to submanifolds of codimension at least two. ∎
The following facts on the maximal connected subgroups of the classical groups can be proven by standard arguments from the representation theory of compact Lie groups, see e.g. . It should be remarked that some of subgroups of $`\mathrm{SO}(n)`$, $`\mathrm{SU}(n)`$ or $`\mathrm{Sp}(n)`$ given by irreducible representations of simple groups of corresponding (real, complex or quaternionic) type are not maximal connected, see for complete lists of inclusions.
###### Proposition 3.6.
Let $`K`$ be a connected proper subgroup of $`\mathrm{SO}(n)`$. Then there is an automorphism $`\alpha `$ of $`\mathrm{SO}(n)`$ such that $`\alpha (K)`$ is contained in one of the following subgroups of $`\mathrm{SO}(n)`$
1. $`\mathrm{SO}(k)\times \mathrm{SO}(nk),`$ $`1kn1`$ .
2. $`\mathrm{SO}(p)\mathrm{SO}(q),`$ $`pq=n,\mathrm{\hspace{0.17em}3}pq`$.
3. $`\mathrm{U}(k),2k=n`$.
4. $`\mathrm{Sp}(p)\mathrm{Sp}(q),`$ $`4pq=n4`$ .
or $`K`$ is a simple irreducible subgroup $`K=\varrho (H)\mathrm{SO}(n)`$, where $`H`$ is a simple compact Lie group and $`\varrho `$ is an irreducible representation of $`H`$ of real type such that $`\mathrm{deg}\varrho =n`$.
###### Proposition 3.7.
Let $`K`$ be a connected proper subgroup of $`\mathrm{SU}(n)`$. Then there is an automorphism $`\alpha `$ of $`\mathrm{SU}(n)`$ such that $`\alpha (K)`$ is contained in one of the following subgroups of $`\mathrm{SU}(n)`$
1. $`\mathrm{SO}(n)`$
2. $`\mathrm{Sp}(m),2m=n`$
3. $`\mathrm{S}(\mathrm{U}(k)\times \mathrm{U}(nk)),1kn1`$
4. $`\mathrm{SU}(p)\mathrm{SU}(q),pq=n,p3,q2`$
or $`K`$ is a simple irreducible subgroup $`K=\varrho (H)\mathrm{SU}(n)`$, where $`H`$ is a simple compact Lie group and $`\varrho `$ is an irreducible representation of $`H`$ of complex type such that $`\mathrm{deg}\varrho =n`$.
###### Proposition 3.8.
Let $`K`$ be a connected proper subgroup of $`\mathrm{Sp}(n)`$. Then there is an automorphism $`\alpha `$ of $`\mathrm{Sp}(n)`$ such that $`\alpha (K)`$ is contained in one of the following subgroups of $`\mathrm{Sp}(n)`$
1. $`\mathrm{U}(n),`$
2. $`\mathrm{Sp}(k)\times \mathrm{Sp}(nk),1kn1`$
3. $`\mathrm{SO}(p)\mathrm{Sp}(q),pq=n,p3,q1`$
or $`K`$ is a simple irreducible subgroup $`K=\varrho (H)\mathrm{Sp}(n)`$, where $`H`$ is a simple compact Lie group and $`\varrho `$ is an irreducible representation of $`H`$ of quaternionic type such that $`\mathrm{deg}\varrho =2n`$.
## 4. Criteria for polarity
The following is a generalization of the criterion for hyperpolarity given in . Note that we do not require the sections to be embedded submanifolds here. Hyperpolar actions are characterized by the property that the Lie triple system $`\nu `$ in Proposition 4.1 is abelian.
###### Proposition 4.1.
Let $`G`$ be a connected compact Lie group, $`KG`$ a symmetric subgroup and let $`\text{g}=\text{k}+\text{p}`$ be the Cartan decomposition. Let $`HG`$ be a closed subgroup. Let $`k`$ be the cohomogeneity of the $`H`$-action on $`G`$. Then the following are equivalent.
1. The $`H`$-action on $`G/K`$ is polar w.r.t some Riemannian metric induced by an $`\mathrm{Ad}(G)`$-invariant scalar product on g.
2. For any $`gG`$ such that $`gK`$ lies in a principal orbit of the $`H`$-action on $`G/K`$ the subspace $`\nu =g^1N_{gK}(HgK)\text{p}`$ is a $`k`$-dimensional Lie triple system such that the Lie algebra $`𝔰=\nu [\nu ,\nu ]`$ generated by $`\nu `$ is orthogonal to $`\mathrm{Ad}(g^1)\text{h}`$.
3. The normal space $`\mathrm{N}_{\mathrm{e}K}(H\mathrm{e}K)\text{p}`$ contains a $`k`$-dimensional Lie triple system $`\nu `$ such that the Lie algebra $`𝔰=\nu [\nu ,\nu ]`$ generated by $`\nu `$ is orthogonal to h.
###### Proof.
Let $`gG`$ be such that $`gK`$ lies in a principal orbit of the $`H`$-action on $`G/K`$. Then the action of $`g^1Hg`$ on $`G/K`$ has a principal orbit containing $`\mathrm{e}K`$ and the equivalence of (i) and (ii) follows from , Proposition, p. 193.
Assume now condition (i) holds. Let $`\mathrm{\Sigma }`$ be a section of the polar $`H`$-action on $`G/K`$ such that $`\mathrm{e}K\mathrm{\Sigma }`$. Let $`\nu =T_{\mathrm{e}K}\mathrm{\Sigma }\text{p}`$, let $`𝔰=\nu [\nu ,\nu ]`$, and let $`S`$ be the connected subgroup of $`G`$ corresponding to $`𝔰`$. Since $`S`$ acts transitively on $`\mathrm{\Sigma }`$, there is an element $`sS`$ such that the point $`sK`$ lies in a principal orbit of the $`H`$-action on $`G/K`$. Now it follows from (ii) that $`\mathrm{Ad}(s^1)\text{h}`$ is orthogonal to the Lie algebra generated by $`s^1N_{sK}(HsK)`$, which coincides with $`𝔰`$. Since $`\mathrm{Ad}(s^1)`$ leaves $`𝔰^{}`$ invariant, we have that h is orthogonal to $`𝔰`$ and (iii) follows.
We will now show that if (iii) holds, then $`\mathrm{\Sigma }=\mathrm{exp}(\nu )G/K`$ meets the orbits orthogonally. Let $`sK\mathrm{\Sigma }`$, where $`s`$ is an arbitrary element of the Lie group $`S`$ corresponding to the Lie algebra generated by $`\nu `$. The tangent space of the $`H`$-orbit through $`sK`$ is orthogonal to $`\mathrm{T}_{sK}\mathrm{\Sigma }`$ if and only if $`s^1\text{h}s\nu `$. But since the adjoint representation of $`G`$ restricted to $`S`$ leaves the orthogonal complement of $`𝔰`$ invariant, $`s^1\text{h}s`$ is perpendicular to $`𝔰`$. Thus the $`H`$-action on $`G/K`$ is polar. ∎
As an immediate consequence of this criterion, the problem of classifying polar actions on $`G/K`$ is reduced to a problem on the Lie algebra level. We conclude this section with the simple observation that a polar action restricted to an invariant totally geodesic submanifold is polar.
###### Lemma 4.2.
Let $`G`$ be a compact Lie group acting polarly on a connected Riemannian manifold $`N`$. Let $`MN`$ be a connected totally geodesic submanifold which is invariant under the $`G`$-action. Then the $`G`$-action on $`M`$ is polar.
###### Proof.
Let $`\mathrm{\Sigma }N`$ be a section of the $`G`$-action on $`N`$. Let $`\mathrm{\Sigma }_0`$ be a connected component of $`\mathrm{\Sigma }M`$. Then the totally geodesic submanifold $`\mathrm{\Sigma }_0M`$ obviously meets the $`G`$-orbits in $`M`$ orthogonally at every intersection point. Furthermore, since $`M`$ is connected, any two orbits of the $`G`$-action on $`M`$ can be joined by a shortest geodesic which meets the principal $`G`$-orbits orthogonally and is hence contained in $`\mathrm{\Sigma }`$ after conjugation with a group element. This geodesic is now also contained in $`M`$, since $`M`$ is totally geodesic. This proves that $`\mathrm{\Sigma }_0`$ meets all $`G`$-orbits in $`M`$. ∎
## 5. Sections and Weyl group actions
Let us first recall some known properties of the Weyl group.
###### Lemma 5.1 (Thorbergsson, Podestà).
Let $`M`$ be a simply connected symmetric space on which a compact, connected Lie group $`G`$ acts polarly and nontrivially. Let $`\mathrm{\Sigma }`$ be a section of the polar action and let $`p\mathrm{\Sigma }`$ be such that the orbit through $`p`$ is singular. Then there is a totally geodesic hypersurface $`H`$ in $`\mathrm{\Sigma }`$ passing through $`p`$ and consisting of singular points; moreover there exists a non-trivial element $`gW_\mathrm{\Sigma }`$ which fixes $`H`$ pointwisely.
The set of singular points in $`\mathrm{\Sigma }`$ is a union of finitely many totally geodesic hypersurfaces $`\{H_i\}_{iI}`$ in $`\mathrm{\Sigma }`$; the Weyl group $`W_\mathrm{\Sigma }`$ is generated by reflections in the hypersurfaces $`\{H_i\}_{iI}`$. Let $`\stackrel{~}{\mathrm{\Sigma }}`$ be the universal covering of $`\mathrm{\Sigma }`$ and let $`\{P_j\}_{jJ}`$ be the collection of all lifts of all the totally geodesic hypersurfaces $`\{H_i\}_{iI}`$ in $`\mathrm{\Sigma }`$. Let $`\widehat{W}_\mathrm{\Sigma }`$ be the subgroup of the isometry group of $`\stackrel{~}{\mathrm{\Sigma }}`$ which is generated by the reflections in the hypersurfaces $`\{P_j\}_{jJ}`$. Then $`\widehat{W}_\mathrm{\Sigma }`$ is a Coxeter group and $`W_\mathrm{\Sigma }`$ is a quotient group of $`\widehat{W}_\mathrm{\Sigma }`$.
###### Proof.
See , Lemma 1A.4 or , Section 2.3 for a more general statement. ∎
The following splitting theorem is a generalization of Lemma 1A.2 in , where $`\stackrel{~}{\mathrm{\Sigma }}_1`$ is a point and the hypothesis is equivalent to a trivial Weyl group action. We consider the weaker hypothesis that the section of a polar action is locally a product such that the Weyl group acts trivially on one factor.
###### Splitting Theorem 5.2.
Let $`N`$ be a compact connected Riemannian symmetric space on which a connected compact Lie group $`G`$ acts polarly. Assume the universal covering $`\stackrel{~}{\mathrm{\Sigma }}`$ of a section $`\mathrm{\Sigma }`$ decomposes as a Riemannian product $`\stackrel{~}{\mathrm{\Sigma }}=\stackrel{~}{\mathrm{\Sigma }}_1\times \stackrel{~}{\mathrm{\Sigma }}_2`$ and the action of $`\widehat{W}_\mathrm{\Sigma }`$ on $`\stackrel{~}{\mathrm{\Sigma }}`$ descends to an action on $`\stackrel{~}{\mathrm{\Sigma }}_1`$ such that
$$w(p,q)=(wp,q)\text{for all }w\widehat{W}_\mathrm{\Sigma }\text{}p\stackrel{~}{\mathrm{\Sigma }}_1\text{}q\stackrel{~}{\mathrm{\Sigma }}_2\text{.}$$
Then the universal cover of $`N`$ is a Riemannian product isometric to $`\stackrel{~}{M}\times \stackrel{~}{\mathrm{\Sigma }}_2`$, where $`M=G\mathrm{\Sigma }_1`$ and $`\mathrm{\Sigma }_1`$ is the image of $`\stackrel{~}{\mathrm{\Sigma }}_1`$ under the covering map $`\stackrel{~}{\mathrm{\Sigma }}\mathrm{\Sigma }`$.
###### Proof.
Let $`\mathrm{\Sigma }`$ be a section and $`p\mathrm{\Sigma }`$ be an arbitrary point of this section. For $`i=1,\mathrm{\hspace{0.17em}2}`$, let $`\mathrm{\Sigma }_i=\mathrm{\Sigma }_i(p)`$ be the totally geodesic submanifolds of $`\mathrm{\Sigma }`$ corresponding to $`\stackrel{~}{\mathrm{\Sigma }}_i`$ (uniquely determined by Theorem 3.1) such that $`p\mathrm{\Sigma }_i`$.
First we show that the isotropy group $`G_p`$ acts trivially on $`\mathrm{T}_p\mathrm{\Sigma }_2`$. Consider the slice representation of $`G_p`$ on $`V=N_p(Gp)`$ which is polar by Proposition 2.3, with section $`\mathrm{T}_p\mathrm{\Sigma }=\mathrm{T}_p\mathrm{\Sigma }_1\mathrm{T}_p\mathrm{\Sigma }_2`$. Now consider the Weyl group $`W^{}`$ of this polar linear representation; it coincides with $`(W_\mathrm{\Sigma })_p`$. Its representation space decomposes into a sum of irreducible modules and one trivial module and the section $`\mathrm{T}_p\mathrm{\Sigma }`$ decomposes accordingly, see . It follows from the hypothesis that $`W^{}`$ acts trivially on the linear subspace $`\mathrm{T}_p\mathrm{\Sigma }_2`$. Since irreducible polar representations have irreducible Weyl groups it follows that $`G_p`$ acts trivially on $`\mathrm{T}_p\mathrm{\Sigma }_2`$.
We will now show that the set $`M(p)=G\mathrm{\Sigma }_1`$ is an embedded submanifold of $`N`$. By the Slice Theorem 2.2 there is an equivariant diffeomorphism $`\mathrm{\Psi }`$ of a $`G`$-invariant open neighborhood around the zero section in the normal bundle $`G\times _{G_p}VG/G_p`$ onto a $`G`$-invariant open neighborhood around the orbit $`Gp`$ such that the zero section in $`G\times _{G_p}V`$ is mapped to the orbit $`Gp`$. The diffeomorphism $`\mathrm{\Psi }`$ is given by the end point map which maps any normal vector $`v_qV=\mathrm{N}_q(Gp)`$ to its image under the exponential map $`\mathrm{exp}_q(v_q)`$.
Since $`\mathrm{\Sigma }_2`$ is a totally geodesic submanifold of $`N`$, we have $`\mathrm{exp}_p(\mathrm{T}_p\mathrm{\Sigma }_2)=\mathrm{\Sigma }_2`$. The subspace $`\mathrm{T}_p\mathrm{\Sigma }_1V`$ is fixed by the Weyl group $`W^{}`$ of the slice representation and hence fixed by $`G_p`$. Hence the orthogonal complement $`S`$ of $`\mathrm{T}_p\mathrm{\Sigma }_2`$ in $`V`$ is a linear subspace invariant under the polar representation of $`G_p`$ on $`V`$. Therefore, $`S`$ defines a smooth subbundle of the normal bundle $`G\times _{G_p}V`$.
But since $`\mathrm{T}_p\mathrm{\Sigma }_1`$ is the section of the $`G_p`$-representation on $`S`$ we have $`S=G_p\mathrm{T}_p\mathrm{\Sigma }_1`$. From the fact that $`\mathrm{\Psi }`$ is an equivariant diffeomorphism it follows now that the elements of the subbundle defined by $`S`$ are mapped into the set $`G\mathrm{\Sigma }_1`$. This shows that, in a neighborhood of $`p`$, the subset $`M(p)=G\mathrm{\Sigma }_1N`$ is a smooth submanifold of codimension $`dim(\mathrm{\Sigma }_2)`$. Thus we see that the symmetric space $`N`$ is foliated by the totally geodesic submanifolds $`\{g\mathrm{\Sigma }_2(p)\}_{gG,p\mathrm{\Sigma }}`$ with integrable normal bundle whose integral manifolds are given by $`\{M(p)\}_{p\mathrm{\Sigma }}`$.
It follows from Theorem A of that the universal cover $`\stackrel{~}{N}`$ of $`N`$ is topologically a product diffeomorphic to $`\stackrel{~}{M}\times \stackrel{~}{\mathrm{\Sigma }}_2`$ such that the projection of $`\stackrel{~}{N}`$ on the factor $`\stackrel{~}{\mathrm{\Sigma }}_2`$ is a Riemannian submersion. We have just shown that the horizontal distribution of this Riemannian submersion is integrable. Since the sectional curvature of $`N`$ is nonnegative, it follows from Theorem 1.3 of that the fibers of this Riemannian submersion are totally geodesic. We conclude that $`\stackrel{~}{N}`$ is a Riemannian product isometric to $`\stackrel{~}{M}\times \stackrel{~}{\mathrm{\Sigma }}_2`$. ∎
###### Corollary 5.3.
Let $`N`$ be an irreducible Riemannian symmetric space of compact type on which a compact Lie group $`G`$ acts polarly and nontrivially. Then the $`G`$-action on $`N`$ has a singular orbit.
###### Proof.
Assume there is no singular orbit. Then by Lemma 5.1 the Weyl group $`W_\mathrm{\Sigma }`$ acts trivially on $`\mathrm{\Sigma }`$. Hence it follows from the Splitting Theorem 5.2 that $`\stackrel{~}{N}`$ is a Riemannian product $`\stackrel{~}{M}\times \stackrel{~}{\mathrm{\Sigma }}`$, where $`M`$ is a $`G`$-orbit. But this is a contradiction to the irreducibility of $`N`$. ∎
The following theorem is a generalization of Proposition 1B.1 of , where it was proved that the section of a polar action on a compact rank one symmetric space has constant curvature.
###### Theorem 5.4.
Let $`N`$ be an irreducible compact simply connected symmetric space on which a compact Lie group $`G`$ acts polarly and non-trivially with section $`\mathrm{\Sigma }`$. Then $`\mathrm{\Sigma }`$ is covered by a Riemannian product of spaces which have constant curvature.
###### Proof.
Let $`\stackrel{~}{\mathrm{\Sigma }}=\stackrel{~}{\mathrm{\Sigma }}_1\times \stackrel{~}{\mathrm{\Sigma }}_2`$ be a decomposition of the universal covering $`\stackrel{~}{\mathrm{\Sigma }}`$ of $`\mathrm{\Sigma }`$ such that $`\mathrm{\Sigma }_1`$ is a Riemannian product of spaces of constant curvature and $`\mathrm{\Sigma }_2`$ is either a point or a Riemannian product of irreducible symmetric spaces of non-constant curvature. The section $`\mathrm{\Sigma }`$ contains a union of finitely many totally geodesic hypersurfaces $`\{H_i\}_{iI}`$ such that the Weyl group $`W_\mathrm{\Sigma }`$ is generated by the reflections in the hypersurfaces $`\{H_i\}_{iI}`$. In view of Corollary 3.5 and since it is well known that the only irreducible symmetric spaces containing totally geodesic hypersurfaces are those of constant curvature, it is clear that the hypothesis of the Splitting Theorem 5.2 is fulfilled and we conclude that $`\mathrm{\Sigma }_2`$ is a point. ∎
It follows from Theorem 5.4 and Lemma 3.3 that the cohomogeneity of a polar action on an irreducible symmetric space $`G/K`$ is less or equal $`\mathrm{rk}(G)+\mathrm{rk}(K)`$. For hermitian symmetric spaces $`G/K`$ the upper bound on the cohomogeneity can be further improved, see Proposition 5.5 below. These dimension bounds are essential for our classification, since they reduce the classification problem to a finite number of cases.
###### Proposition 5.5.
Let $`H`$ be a compact Lie group acting polarly on a compact Kähler manifold $`M`$. Then the cohomogeneity of the $`H`$-action on $`M`$ is less or equal $`\mathrm{rk}(H)`$.
###### Proof.
By the Equivalence Theorem , see also Theorem 1.4 in , the cohomogeneity of the $`H`$-action is equal to the difference between the rank of $`H`$ and the rank of a regular isotropy subgroup of $`H`$. ∎
We have the following lower bounds on the dimension of groups acting polarly on the classical symmetric spaces.
###### Proposition 5.6.
Let $`H`$ be a connected compact Lie group acting polarly and non-trivially on a symmetric space $`M`$. Assume $`3kn3`$, $`2\mathrm{}n2`$ and let $`d=dim(H)`$.
1. If $`M=𝔾_k(^n)`$, then $`d2n9`$.
2. If $`M=𝔾_{\mathrm{}}(^n)`$, then $`d3n7`$.
3. If $`M=𝔾_{\mathrm{}}(^n)`$, then $`d6n16`$.
4. If $`M=\mathrm{SO}(n)/\mathrm{U}(\frac{n}{2})`$, then $`d\frac{n^2}{4}n`$.
5. If $`M=\mathrm{SU}(n)/\mathrm{SO}(n)`$, then $`d\frac{n^2}{2}n`$.
6. If $`M=\mathrm{SU}(n)/\mathrm{Sp}(\frac{n}{2})`$, then $`d\frac{n^2}{2}2n`$.
7. If $`M=\mathrm{Sp}(n)/\mathrm{U}(n)`$, then $`dn^2`$.
###### Proof.
Follows from Proposition 5.5 in case of the spaces $`𝔾_{\mathrm{}}(^n)`$, $`\mathrm{SO}(n)/\mathrm{U}(\frac{n}{2})`$, and $`\mathrm{Sp}(n)/\mathrm{U}(n)`$, which are Hermitian symmetric, and from Theorem 5.4 and Lemma 3.3 otherwise. ∎
## 6. Polar subactions
In this section, we will introduce our main tool for classifying polar actions through studying slice representations. The basic observation is the following maximality property of linear polar actions, see , Theorem 6.
###### Theorem 6.1.
Let $`G\mathrm{SO}(n)`$ be a closed connected subgroup which acts irreducibly on $`^n`$ and non-transitively on the sphere $`\mathrm{S}^{n1}^n`$. Let $`HG`$ be a closed connected subgroup $`\{\mathrm{e}\}`$ that acts polarly on $`^n`$. Then the $`H`$-action and the $`G`$-action on $`^n`$ are orbit equivalent.
The proof of the above theorem relies on . As an immediate consequence of Theorem 6.1 we have the following, cf. , Theorem 2.2.
###### Corollary 6.2.
Let $`X`$ be a strongly isotropy irreducible Riemannian homogeneous space. Assume a connected compact Lie group $`H`$ acts polarly on $`X`$ such that the $`H`$-action has a one-dimensional orbit $`Hp`$ or a fixed point $`pX`$. Then the space $`X`$ is locally symmetric. Furthermore, $`X`$ is a rank-one symmetric space or the action of $`H`$ is orbit equivalent to the action of the connected component of the isotropy group of $`X`$ at $`p`$.
###### Proof.
Assume first that $`pX`$ is a fixed point of the $`H`$-action on $`X`$. Let $`K`$ be the connected component of the isotropy group of $`p`$. The isotropy representation of $`K`$ on $`\mathrm{T}_pX`$ restricted to $`H`$ is polar by Proposition 2.3. If the action of $`K`$ on the unit sphere in $`\mathrm{T}_pX`$ is transitive, then the space $`X`$ is rank-one symmetric. If $`K`$ does not act transitively on the sphere, then the linear $`H`$-action on $`\mathrm{T}_pX`$ is orbit equivalent to the $`K`$-action by Theorem 6.1, in particular they have the same cohomogeneity; hence the $`K`$-action on $`\mathrm{T}_pX`$ is polar. It now follows that the principal orbits of the $`H`$-action agree with those of the $`K`$-action on $`X`$ and the orbit equivalence of the two actions follows from Proposition 4.1, since the principal orbits of a hyperpolar action determine all other orbits. In case $`X`$ is compact, the symmetry follows from , since then one may assume that $`X`$ is a homogeneous space of a simple compact Lie group, see , Chapter I.1. Non-compact strictly isotropy irreducible Riemannian homogeneous spaces are symmetric by .
Now assume $`pX`$ is such that $`dim(Hp)=1`$. If $`Hp`$ is a regular orbit, it follows that a section $`\mathrm{\Sigma }X`$ is a totally geodesic hypersurface and hence $`X`$ is locally isometric to a space of constant curvature. Assume now that $`Hp`$ is a singular orbit, hence the slice representation of $`H_p`$ on $`\mathrm{N}_p(Hp)`$ is nontrivial and polar by Proposition 2.3. However, since $`\mathrm{T}_p(Hp)`$ is one-dimensional, the isotropy representation of $`H_p`$ on $`\mathrm{T}_pX=\mathrm{T}_p(Hp)\mathrm{N}_p(Hp)`$ is polar. It now follows from Theorem 6.1 that the irreducible isotropy representation of $`M`$ at $`p`$ is orbit equivalent to the reducible $`H_p`$-action on $`\mathrm{T}_pM`$, a contradiction. ∎
In particular, we may restrict our attention to actions without fixed point in the following. The full classification of connected Lie groups acting polarly with a fixed point on the irreducible symmetric spaces of higher rank follows immediately from Lemma 2.6. As the proof of Corollary 6.2 shows, one obtains the same result also under the weaker hypothesis that the linear action of $`H`$ on the tangent space $`\mathrm{T}_pX`$ is polar.
Let $`G`$ be a connected compact Lie group acting isometrically on a Riemannian manifold. We say the action of $`G`$ on $`M`$ is polarity minimal if there is no closed connected subgroup $`HG`$ which acts nontrivially and polarly on $`M`$ and such that the $`H`$-action is not orbit equivalent to the $`G`$-action. Note that a polarity minimal action can be polar or non-polar. We give various sufficient conditions for an orthogonal representation to be polarity minimal in the following proposition.
###### Proposition 6.3.
Let $`\rho :G\mathrm{O}(V)`$ be a representation of the compact connected Lie group $`G`$. Then $`\rho `$ is polarity minimal if one of the following holds.
1. The representation $`\rho `$ is irreducible of cohomogeneity $`2`$.
2. The representation space $`V`$ is the direct sum of two equivalent $`G`$-modules.
3. The representation space $`V`$ contains a $`G`$-invariant submodule $`W`$ such that the $`G`$-representation on $`W`$ is almost effective, non-polar, and polarity minimal.
###### Proof.
Part (i) is a just a reformulation of Theorem 6.1. Assume now $`V`$ is the direct sum of two equivalent $`G`$-modules; then the representation $`\rho `$ restricted to any closed connected subgroup $`HG`$ which acts nontrivially on $`V`$ will have two equivalent nontrivial submodules; it then follows from , Lemma 2.9 that $`H`$ acts non-polarly on $`V`$; this proves part (ii). To prove part (iii), assume there is a closed connected subgroup $`H`$ of $`G`$ acting polarly on $`V`$; since the $`G`$-action on the subspace $`W`$ is non-polar and polarity minimal, it follows that $`H`$ acts trivially on $`W`$. But $`W`$ is an almost effective representation, thus $`H`$ acts trivially on all of $`V`$. ∎
While we do not have an a priori proof that polar actions on irreducible compact symmetric spaces of higher rank are orbit maximal, the following proposition gives various sufficient conditions under which one can show that certain non-polar actions are polarity minimal. In fact, this is our main tool to exclude subactions and it will be used frequently in the sequel.
###### Lemma 6.4.
Let $`G`$ be compact Lie group and $`KG`$ be symmetric subgroup such that $`M=G/K`$ is an irreducible symmetric space and let $`HG`$ be a closed subgroup. The action of $`H`$ on $`M`$ is non-polar and polarity minimal if there is a non-polar polarity minimal submodule $`V\mathrm{N}_p(Hp)`$ of the slice representation at $`p`$ such that one of the following holds.
1. $`M`$ is Hermitian symmetric and $`dim(V)>\mathrm{rk}(H)`$.
2. $`dim(V)>s(M)`$, where $`s(M)`$ is the maximal dimension of a totally geodesic submanifold of $`M`$ locally isometric to a product of spaces with constant curvature, cf. Lemma 3.3.
3. $`V\text{p}=\mathrm{T}_pM`$ (where $`\text{g}=\text{k}\text{p}`$ as usual such that k is the Lie algebra of $`K=G_p`$) contains a Lie triple system corresponding to an irreducible symmetric space of nonconstant curvature, e.g. an irreducible symmetric space of higher rank.
4. The isotropy group $`HK`$ acts almost effectively on $`V`$ and $`\mathrm{rk}(HK)=\mathrm{rk}(H)`$.
###### Proof.
Assume a closed connected subgroup $`UH`$ acts polarly on $`M`$. Consider the isotropy group at $`U_p`$ of the $`U`$-action on $`M`$. Since $`U_pH_p`$, the action of $`U_p`$ on the normal space $`\mathrm{N}_p(Up)`$ leaves the subspace $`V`$ invariant. By Proposition 2.3, the slice representation of $`U_p`$ on $`\mathrm{N}_p(Up)`$ is polar, in particular, the $`U_p`$-action on $`V`$ is polar. Since the action of $`H_p`$ on $`V`$ is polarity minimal and non-polar, it follows that the action of the connected component $`\left(U_p\right)_0`$ on $`V`$ is trivial. Hence $`V`$ is contained in the section of the polar $`U_p`$-action on $`\mathrm{N}_p(Up)`$ and thus $`V`$ is tangent to a section $`\mathrm{\Sigma }`$ of the $`U`$-action on $`M`$; in particular, $`dim(\mathrm{\Sigma })dim(V)`$. Part (i) now follows from Lemma 5.5. Parts (ii) and (iii) follow from Theorem 5.4. If $`\mathrm{rk}(HK)=\mathrm{rk}(H)`$ and $`HK`$ acts almost effectively on $`V`$, then any closed subgroup $`UH`$ with $`dimU>0`$ will have an intersection $`UKHK`$ of positive dimension with $`K`$; but $`UK`$ acts on $`V`$ non-polarly since $`V`$ is polarity minimal; this proves (iv). ∎
## 7. Hermann actions of higher cohomogeneity
In the remaining part of the paper, we will carry out the classification. We begin with subactions of Hermann actions whose cohomogeneity is $`2`$.
To study actions of reducible groups on the Grassmannians we will need the following technical lemma. Let us first introduce some notation. Let $`G\mathrm{G}\mathrm{}(n,)`$. Let $`V`$ be a linear subspace of $`^n`$. Then we define the normalizer of $`V`$ in $`G`$ as $`N_G(V)=\{gGg(V)=V\}`$, and similarly by $`Z_G(V)=\{gGg|_V=\mathrm{id}_V\}`$ the centralizer of $`V`$ in $`G`$. Clearly, $`N_G(V)`$ is a subgroup of $`G`$ and since the elements of $`N_G(V)`$ leave $`V`$ invariant, the group $`N_G(V)`$ acts on $`V`$. The kernel of this representation is the normal subgroup $`Z_G(V)N_G(V)`$. The group $`N_G(V)`$ is the isotropy subgroup $`G_V`$ of the $`G`$-action on the Grassmannian $`𝔾_{dimV}(^n)`$ of $`(dimV)`$-dimensional linear subspaces in $`^n`$.
###### Lemma 7.1.
Let $`H\mathrm{SO}(n)`$ be a closed connected proper subgroup.
1. If for any $`8`$-dimensional subspace $`V^n`$ the natural action of the connected component of $`N_H(V)/Z_H(V)`$ on $`V`$ is equivalent to the 8-dimensional spin representation of $`\mathrm{Spin}(7)`$ or the standard representation of $`\mathrm{SO}(8)`$ then $`n=8`$ and $`H\mathrm{Spin}(7)`$.
2. If for any $`7`$-dimensional subspace $`V^n`$ the natural action of the connected component of $`N_H(V)/Z_H(V)`$ on $`V`$ is equivalent to the 7-dimensional irreducible representation of $`\mathrm{G}_2`$ or the standard representation of $`\mathrm{SO}(7)`$ then either $`n=7`$ and $`H\mathrm{G}_2`$ or $`n=8`$ and $`H\mathrm{Spin}(7)`$.
###### Proof.
We first show that in both cases the group $`H`$ acts transitively on the unit sphere in $`^n`$. Let $`p`$, $`q^n`$ be two unit vectors and let $`V`$ be linear subspace of $`^n`$ containing $`p`$ and $`q`$ such that $`V`$ is $`8`$\- or $`7`$-dimensional, respectively. Then it follows from the hypothesis that $`N_H(V)/Z_H(V)`$ acts transitively on the unit sphere in the space $`V`$, thus there is an element in $`H`$ which maps $`p`$ to $`q`$. This shows that $`H`$ acts transitively on the unit sphere in $`^n`$ and hence the pair $`(H,n)`$ is one of the following, see Table 7 in .
| $`H`$ | $`\mathrm{U}(m),`$ | $`\mathrm{SU}(m),`$ | $`\mathrm{Sp}(\mathrm{})\mathrm{Sp}(1),`$ | $`\mathrm{Sp}(\mathrm{})\mathrm{U}(1),`$ | $`\mathrm{Sp}(\mathrm{}),`$ | $`\mathrm{Spin}(7)`$ | $`\mathrm{Spin}(9)`$ | $`\mathrm{G}_2`$ |
| --- | --- | --- | --- | --- | --- | --- | --- | --- |
| | $`m2`$ | $`m2`$ | $`\mathrm{}2`$ | $`\mathrm{}2`$ | $`\mathrm{}2`$ | | | |
| $`n`$ | $`2m`$ | $`2m`$ | $`4\mathrm{}`$ | $`4\mathrm{}`$ | $`4\mathrm{}`$ | $`8`$ | $`16`$ | $`7`$ |
It is easy to see that the first five groups do not have the property described in the hypothesis. For the groups $`\mathrm{Spin}(7)`$ and $`\mathrm{G}_2`$ the statement is either trivial or follows from the well-known fact $`\mathrm{S}^7=\mathrm{Spin}(7)/\mathrm{G}_2`$.
It remains the case of $`H=\mathrm{Spin}(9)`$ acting on $`^{16}`$. To prove the assertion for $`dimV=8`$, it suffices to exhibit an isotropy subgroup of the $`\mathrm{Spin}(9)`$-action on $`𝔾_8(^{16})`$ not containing $`\mathrm{Spin}(7)`$ as a Lie subgroup. First choose an $`8`$-dimensional subspace $`\stackrel{~}{V}^{16}`$ such that the subgroup $`\mathrm{Spin}(8)`$, acting by a representation equivalent to the sum of the two half-spin representations on $`^{16}`$, stabilizes $`\stackrel{~}{V}`$. Since $`\mathrm{Spin}(8)\mathrm{Spin}(9)`$ is maximal connected, it coincides with the connected component of the isotropy group $`H_{\stackrel{~}{V}}`$ of the $`H`$-action on the Grassmannian $`𝔾_8(^{16})`$. Thus the $`H`$-orbit through $`\stackrel{~}{V}`$ is $`8`$-dimensional. We will determine the slice representation of the $`H`$-action at $`\stackrel{~}{V}`$. The group $`\left(H_{\stackrel{~}{V}}\right)_0\mathrm{Spin}(8)`$ acts on the tangent space $`\mathrm{T}_{\stackrel{~}{V}}𝔾_8(^{16})`$ by the tensor product of the two half-spin representations of $`\mathrm{Spin}(8)`$. By Weyl’s dimension formula, this representation contains an irreducible summand which is $`56`$-dimensional and must therefore coincide with the normal space $`\mathrm{N}_{\stackrel{~}{V}}(H\stackrel{~}{V})`$. This shows that $`\stackrel{~}{V}`$ lies in a singular orbit of the $`H`$-action on $`𝔾_8(^{16})`$. By , the principal isotropy subgroups of this slice representation are finite and we conclude that for generic subspaces $`V^{16}`$ the group $`N_H(V)/Z_H(V)`$ is finite.
Similarly, to prove the assertion for $`H=\mathrm{Spin}(9)`$ and $`dimV=7`$, choose a $`7`$-dimensional subspace $`\stackrel{~}{V}^{16}`$ such that $`\mathrm{Spin}(7)\mathrm{Spin}(9)`$ stabilizes $`\stackrel{~}{V}`$. Using an analogous argument as in the case $`k=8`$ we see that the $`48`$-dimensional slice representation at $`\stackrel{~}{V}`$ of the $`H`$-action on $`𝔾_7(^{16})`$ has finite principal isotropy subgroups and hence for a generic $`7`$-dimensional subspace $`V^{16}`$ the group $`N_H(V)/Z_H(V)`$ is finite. ∎
###### Lemma 7.2.
Let $`H`$, $`G`$, $`K`$ be as in the following table, where $`2k,\mathrm{}\frac{n}{2}`$.
| $`H`$ | $`G/K`$ |
| --- | --- |
| $`\mathrm{S}(\mathrm{U}(k)\times \mathrm{U}(nk))`$, | $`\mathrm{SU}(n)/\mathrm{S}(\mathrm{U}(\mathrm{})\times \mathrm{U}(n\mathrm{}))`$ |
| $`\mathrm{SO}(k)\times \mathrm{SO}(nk)`$ | $`\mathrm{SO}(n)/\mathrm{SO}(\mathrm{})\times \mathrm{SO}(n\mathrm{})`$ |
| $`\mathrm{Sp}(k)\times \mathrm{Sp}(nk)`$ | $`\mathrm{Sp}(n)/\mathrm{Sp}(\mathrm{})\times \mathrm{Sp}(n\mathrm{})`$ |
Let $`U`$ be a connected subgroup of $`H`$. Then the action of $`U`$ on $`G/K`$ is polar if either $`U`$ = $`H`$ or $`U`$ is conjugate to one of the following subgroups, where in each case the $`U`$-action on $`G/K`$ is orbit equivalent to the $`H`$-action. In particular, the $`U`$-action on $`G/K`$ is hyperpolar.
| $`U`$ | $`G`$ | Range |
| --- | --- | --- |
| $`\mathrm{G}_2\times \mathrm{G}_2`$ | $`\mathrm{SO}(14)`$ | $`\mathrm{}=2`$ |
| $`\mathrm{G}_2\times \mathrm{Spin}(7)`$ | $`\mathrm{SO}(15)`$ | $`\mathrm{}=2`$ |
| $`\mathrm{Spin}(7)\times \mathrm{Spin}(7)`$ | $`\mathrm{SO}(16)`$ | $`\mathrm{}=2,3`$ |
| $`\mathrm{G}_2\times \mathrm{SO}(n7)`$ | $`\mathrm{SO}(n)`$ | $`\mathrm{}=2,n9`$ |
| $`\mathrm{Spin}(7)\times \mathrm{SO}(n8)`$ | $`\mathrm{SO}(n)`$ | $`\mathrm{}=2,3,n10`$ |
| $`\mathrm{SU}(k)\times \mathrm{SU}(nk)`$ | $`\mathrm{SU}(n)`$ | $`(k,\mathrm{})(\frac{n}{2},\frac{n}{2})`$ |
###### Proof.
To prove the lemma, we compute certain slice representations, cf. Section 2.3 in . We assume in the following that the maximal reducible groups are standardly embedded as block diagonal matrices, cf. (2.1). Assume first $`k\mathrm{}`$. We compute a slice representation of the action of $`H=\mathrm{SO}(k)\times \mathrm{SO}(nk)`$ on $`G/K=\mathrm{SO}(n)/\mathrm{SO}(\mathrm{})\times \mathrm{SO}(n\mathrm{})`$. The connected component of the isotropy group is the group $`(HK)_0=\mathrm{SO}(k)\times \mathrm{SO}(\mathrm{}k)\times \mathrm{SO}(n\mathrm{})`$; it acts on the normal space
$$N_{\mathrm{e}K}(H\mathrm{e}K)=\left\{\left(\begin{array}{ccc}0& 0& M\\ 0& 0& 0\\ M^t& 0& 0\end{array}\right)\right|M^{k\times n\mathrm{}}\}𝔰𝔬(n)$$
by the tensor product of the two standard representations of the first and the last factor i.e. $`\mathrm{SO}(k)\mathrm{SO}(n\mathrm{})`$. Assume $`UH`$ is a closed subgroup acting polarly on $`G/K`$. By Lemma 2.6 the connected component of the isotropy group $`UK`$ of the $`U`$-action must contain the product of the first and last factor $`L=\mathrm{SO}(k)\times \mathrm{SO}(n\mathrm{})HK`$, except possibly in cases $`k=2,\mathrm{\hspace{0.17em}3}`$ and $`n\mathrm{}=7,\mathrm{\hspace{0.17em}8}`$ which will be treated below. Since this argument also holds for any conjugate subgroup $`hUh^1`$, $`hH`$, it follows that $`U`$ contains $`hLh^1`$, for all $`hH`$, hence $`U`$ contains the subgroup generated by $`\{hxh^1hH,xL\}`$, which is the minimal normal subgroup of $`H`$ containing $`L`$ and we conclude $`H=U`$.
Using an analogous argument in the case of $`H=\mathrm{S}(\mathrm{U}(k)\times \mathrm{U}(nk))`$ acting on $`G/K=\mathrm{SU}(n)/\mathrm{S}(\mathrm{U}(\mathrm{})\times \mathrm{U}(n\mathrm{}))`$ we see that the only polar subaction is the action of $`U=\mathrm{SU}(k)\times \mathrm{SU}(nk)`$ except in the case of $`U=\mathrm{SU}(k)\times \mathrm{SU}(k)`$ acting on $`G/K=\mathrm{SU}(2k)/\mathrm{S}(\mathrm{U}(k)\times \mathrm{U}(k))`$, where the slice representation of the $`U`$-action is non-polar, see Lemma 2.6. The same argument also works for the actions on the quaternionic Grassmannians and for the case $`\mathrm{}k`$.
It remains to study the case where a slice representation of the $`U`$-action is given by the first three rows of Lemma 2.6. It follows from Lemma 7.1 that this can only happen if $`U`$ is obtained from $`H`$ by replacing an $`\mathrm{SO}(7)`$-factor with $`\mathrm{G}_2`$ or replacing an $`\mathrm{SO}(8)`$-factor with $`\mathrm{Spin}(7)`$. A dimension count shows that the actions obtained in this fashion are orbit equivalent to the respective $`H`$-action. ∎
###### Theorem 7.3.
Let $`G`$ be a connected simple compact Lie group and let $`H`$ and $`K`$ be two non-conjugate connected symmetric subgroups of $`G`$ such that the cohomogeneity $`r`$ of the $`H`$-action on $`G/K`$ is $`2`$. Let $`UH`$ be a closed connected nontrivial subgroup acting polarly on $`G/K`$. Then the action of $`U`$ on $`G/K`$ is orbit equivalent to the hyperpolar $`H`$-action on $`G/K`$.
Furthermore, $`UH`$ if and only if $`U`$ is as described in Lemma 7.2 or the triple $`(U,G,K)`$ is one of $`(\mathrm{SU}(2n2k1)\times \mathrm{SU}(2k+1),\mathrm{SU}(2n),\mathrm{Sp}(n))`$; $`(\mathrm{SU}(n),\mathrm{Sp}(n),`$ $`\mathrm{Sp}(k)`$ $`\times \mathrm{Sp}(nk))`$; $`(\mathrm{SU}(n),\mathrm{SO}(2n),\mathrm{SO}(k)\times \mathrm{SO}(2nk))`$, $`k<n`$; $`(\mathrm{SU}(n),\mathrm{SO}(2n),`$ $`\alpha (\mathrm{U}(n)))`$; $`(\mathrm{Spin}(10),\mathrm{E}_6,\mathrm{SU}(6)\mathrm{Sp}(1))`$; $`(\mathrm{E}_6,\mathrm{E}_7,\mathrm{SO}^{}(12)\mathrm{Sp}(1))`$.
###### Proof.
To prove the theorem, we use the explicit knowledge of slice representations of Hermann actions as given in Tables 3 and 5. The first table is a list of all Hermann actions (i.e. a list of all pairs $`(H,K)`$ of non-conjugate symmetric subgroups of the simple compact Lie groups $`G`$ up to automorphisms of $`G`$); Table 5 contains information about one irreducible slice representation of each action; slice representations of Hermann actions are s-representations by Lemma 11.1 and so each representation is described by a symmetric space $`G^{}/K^{}`$ whose isotropy representation $`\chi (G^{},K^{})`$ is equivalent to the slice representation on the Lie algebra level; in the third column, the (local isomorphism type of the) kernel of the slice representation is given. It is straightforward to determine these slice representations for actions on the classical symmetric spaces, for the exceptional symmetric spaces one may use the technique described in Remark 10.1, cf. also , Prop. 3.5. Note that actually two different actions are given in each row of the table, i.e. the action of $`H`$ on $`G/K`$ and the action of $`K`$ on $`G/H`$; they have the same isotropy subgroups and slice representations.
Assume now that $`H`$ and $`K`$ are symmetric subgroups of the simple compact Lie group $`G`$ and $`UH`$ is a closed connected subgroup acting polarly on $`G/K`$ and such that the hyperpolar action of $`H`$ on $`G/K`$ is of cohomogeneity $`r2`$. The subactions of the types A III-III, BD I-I, and C II-II were treated in Lemma 7.2.
The slice representations given by the table are irreducible and non-transitive on the sphere, since we assume the cohomogeneity is $`2`$, thus we may apply Theorem 6.1. The isotropy group of the $`U`$-action at $`p=\mathrm{e}KG/K`$ is $`U_p=UH_p`$. The representation of $`H_p=HK`$ on the normal space $`N_p(Hp)`$ restricted to $`U_p`$ occurs as a submodule in the slice representation of $`U_p`$ on $`N_p(Up)`$, and is therefore polar. By Theorem 6.1, the $`U_p`$-action on $`V=N_p(Hp)`$ is either orbit equivalent to the $`H_p`$-action or trivial.
We first show that the slice representation of $`H_p`$ restricted to $`U_p`$ is non-trivial. If it is trivial, then $`V`$ is contained in the tangent space of a section through $`p`$ and we obtain a contradiction with Theorem 5.4 since $`V\text{g}`$ is a Lie triple system corresponding to a totally geodesic submanifold of $`G/K`$ which is isometric to an irreducible symmetric space of higher rank, see the proof of Lemma 11.1. We may therefore assume the $`U_p`$-action on $`V`$ is locally orbit equivalent to the irreducible polar representation of $`H_p`$. From Table 5 we see that we may assume the slice representation of $`H_p`$ on $`V`$ is not equivalent to one of the first three items in Lemma 2.6, except in the case of A I-III, which will be treated separately.
Then it follows that the group $`U_p`$ contains the (component of the) isotropy group $`H_p`$ or a subgroup as described in the 4th, 5th, and 6th item of Lemma 2.6. In these cases there exists a uniquely<sup>1</sup><sup>1</sup>1In the first three items of Lemma 2.6, an orbit equivalent subgroup is only unique up to conjugation. determined connected subgroup $`LH_p`$ which is minimal with respect to the property that the $`L`$-action on $`V`$ is orbit equivalent to the $`H_p`$-action. Note that this argument actually shows that for any $`hH`$ also $`hUh^1`$ contains the subgroup $`LH_p`$, hence $`U`$ contains all groups $`h^1Lh`$ conjugate to $`L`$ in $`H`$. We conclude that $`U`$ contains the subgroup $`\widehat{L}`$ generated by $`\{h\mathrm{}h^1\mathrm{}L,hH\}`$, i.e. the minimal normal subgroup of $`H`$ containing $`L`$.
#### Subgroups of codimension one
Let us first consider the case where $`UH`$ is a subgroup of codimension one, i.e. $`H=U\mathrm{U}(1)`$, then we have that either $`U`$ acts transitively on the $`H`$-orbit through $`p`$, in which case the $`U`$-action and the $`H`$-action are orbit equivalent, or $`U`$ acts with cohomogeneity one on the orbit $`Hp`$ in which case we arrive at a contradiction since a section $`\mathrm{\Sigma }`$ through $`p`$ of the $`U`$-action contains the flat section $`\mathrm{\Sigma }_0`$ of the $`H`$-action, on whose tangent space $`\mathrm{T}_p\mathrm{\Sigma }_0`$ the Weyl group of the irreducible slice representation still acts irreducibly when restricted to the Weyl group of the $`U_p`$-representation, so $`\mathrm{\Sigma }`$ would be either flat, contradicting Proposition 2.8, or an irreducible symmetric space of dimension $`r+1`$ and rank $`r2`$, which does not exist.
#### Subactions of Hermann actions on exceptional symmetric spaces
Assume the subgroup $`UH`$ acts polarly on the symmetric space $`G/K`$. One can see from Table 3 that the group $`H`$ has either one or two simple factors if it is semisimple or it is the product $`H=H^{}\mathrm{U}(1)`$ of a one-dimensional abelian and a simple factor. Since $`U`$ contains the nontrivial normal subgroup $`\widehat{L}`$ of $`H`$ it follows that $`U=H`$ if $`H`$ is simple; if $`H=H^{}\mathrm{U}(1)`$ then $`U`$ contains $`H^{}`$ (since $`dim\widehat{L}>1`$) and the $`U`$-action is orbit equivalent to the $`H`$-action by the argument above since then $`UH`$ is a subgroup of codimension one; in those cases where $`H`$ is a product of two simple factors, comparison of the Tables 4 and 5 shows that in each case, except for $`H=\mathrm{SU}(6)\mathrm{SU}(2)`$, $`G/K=\mathrm{E}_6/\mathrm{Spin}(10)\mathrm{U}(1)`$, the normal subgroup $`\widehat{L}`$ contains both simple factors of $`H`$ and it follows that $`H=U`$. Consider the action of $`H=\mathrm{SU}(6)\mathrm{SU}(2)`$ on $`G/K=\mathrm{E}_6/\mathrm{Spin}(10)\mathrm{U}(1)`$; in this case it follows from the data given in Table 5 only that $`\widehat{L}`$ contains the $`\mathrm{SU}(6)`$-factor of $`H`$. An explicit calculation as described in Remark 10.1 shows that the embedding of the connected component of the isotropy group $`(HK)_0=\mathrm{U}(1)\mathrm{SU}(4)\mathrm{SU}(2)\mathrm{SU}(2)`$ into $`H=\mathrm{SU}(6)\mathrm{SU}(2)`$ is such that the $`\mathrm{SU}(2)`$-factor in the kernel of the slice representation lies in the $`\mathrm{SU}(6)`$-factor of $`H`$, and the other $`\mathrm{SU}(2)`$-factor of $`(HK)_0`$, which acts nontrivially on the slice, coincides with the $`\mathrm{SU}(2)`$-factor of $`H`$. From this it follows that the actions $`\mathrm{SU}(6)`$ or $`\mathrm{SU}(6)\mathrm{U}(1)`$ on $`G/K`$ have a slice representation with two equivalent nontrivial modules and are therefore not polar.
#### Subactions of Hermann actions on classical symmetric spaces
The cases A I-II, A II-III, C I-II, D I-III, D III-III’, and D<sub>4</sub> I-I’ can be handled in a similar way as the subactions on the exceptional spaces. One can see from Table 5 that $`\widehat{L}`$ contains every simple factor of $`H`$. For the case of D<sub>4</sub> I-I’, i.e. subactions of $`\mathrm{Spin}(5)\mathrm{Spin}(3)\mathrm{Sp}(2)\mathrm{Sp}(1)`$ on $`\mathrm{SO}(8)/\mathrm{SO}(3)\times \mathrm{SO}(5)`$, the slice representation was explicitly computed in , p. 592-593, and it follows that $`\widehat{L}`$ is not contained in one of the simple factors of $`H`$, thus $`H=U`$.
It remains to study the case A I-III. For the slice representation of this action there are in some cases orbit equivalent polar subgroups as given in the first three items of Lemma 2.6; otherwise the argument is as above. Assume $`H=\mathrm{SO}(n)`$, $`G=\mathrm{SU}(n)`$, $`K=\mathrm{S}(\mathrm{U}(k)\times \mathrm{U}(nk))`$, where $`(n,k)=(9,2)`$, $`(10,2)`$ or $`(11,3)`$. Let us first consider the $`H`$-action on $`G/K`$. The connected component of the isotropy subgroup at $`\mathrm{e}K`$ is $`\mathrm{SO}(k)\times \mathrm{SO}(nk)`$. It follows that $`U`$ must contain the group given in the right column of the table in Lemma 2.6 and it follows from Lemma 7.4 below that either $`U=H`$ or $`UH_p`$, but in the latter case the $`U`$-action on $`G/K`$ has a fixed point. Finally, consider the $`K`$-action on $`G/H`$, i.e. assume a closed connected subgroup $`UK=\mathrm{S}(\mathrm{U}(k)\times \mathrm{U}(nk))`$ acts polarly on $`\mathrm{SU}(n)/\mathrm{SO}(n)`$; it follows by the arguments above that $`U`$ contains a subgroup $`L`$ conjugate to $`\mathrm{SO}(2)\times \mathrm{G}_2`$, if $`K=\mathrm{S}(\mathrm{U}(2)\times \mathrm{U}(7))`$, or $`\mathrm{SO}(k)\times \mathrm{Spin}(7)`$, if $`K=\mathrm{S}(\mathrm{U}(k)\times \mathrm{U}(8))`$; all possibilities for the group $`U`$ are given by Lemma 7.4. It follows that the slice representation $`V|_{U_p}\chi (K,U)|_{U_p}`$ of the $`U`$-action on $`G/H`$ is non-polar by if $`U`$ does not contain both simple factors of $`K`$, thus the codimension of $`U`$ in $`K`$ is at most one and we conclude that the $`U`$-action is orbit equivalent to the $`K`$-action.
To prove the last part of the theorem, one can easily determine all proper closed subgroups $`U`$ of $`H`$ whose action on $`G/K`$ is orbit equivalent to the $`H`$-action on $`G/K`$ by using the information from Table 5. ∎
For the proof of Theorem 7.3 we used the following simple Lemma.
###### Lemma 7.4.
For the following inclusions of compact Lie groups $`ABC`$, the intermediate subgroups $`B`$ are unique in the following sense: If $`B^{}C`$ is a closed connected subgroup such that $`AB^{}C`$ then $`B^{}=B`$.
$`\mathrm{G}_2\mathrm{SO}(7)`$ $`\mathrm{SU}(7);`$
$`\mathrm{Spin}(7)\mathrm{SO}(8)`$ $`\mathrm{SU}(8);`$
$`\mathrm{SO}(k)\times \mathrm{G}_2\mathrm{SO}(k)\times \mathrm{SO}(7)`$ $`\mathrm{SO}(7+k),k;`$
$`\mathrm{SO}(k)\times \mathrm{Spin}(7)\mathrm{SO}(k)\times \mathrm{SO}(8)`$ $`\mathrm{SO}(8+k),k.`$
###### Proof.
It is easily checked in each case that the representation $`\chi (C,A)`$ splits into the irreducible modules $`\chi (B,A)`$ and $`\chi (C,B)|_A`$. Note that $`\chi (\mathrm{SU}(7),\mathrm{SO}(7))|_{\mathrm{G}_2}`$ is equivalent to the irreducible $`27`$-dimensional representation of $`\mathrm{G}_2`$ and $`\chi (\mathrm{SU}(8),\mathrm{SO}(8))|_{\mathrm{Spin}(7)}`$ is equivalent to the irreducible $`35`$-dimensional representation of $`\mathrm{Spin}(7)`$, cf. Table 1, p. 364 of . ∎
## 8. Actions of non-simple irreducible groups
In the following, we will assume that $`G`$ is a simple classical compact Lie group $`G=\mathrm{SO}(n)`$, $`\mathrm{SU}(n)`$ or $`\mathrm{Sp}(n)`$ and $`K`$ is a symmetric subgroup such that $`\mathrm{rk}(G/K)2`$. We will classify all closed connected subgroups $`HG`$ such that $`H`$ acts polarly on $`G/K`$.
The symmetric quotient spaces of the simple classical compact Lie groups of rank $`2`$ which are not locally isometric to one of the Grassmannians $`𝔾_k(^n)`$, $`𝔾_k(^n)`$, $`𝔾_k(^n)`$ are the following:
(8.1)
$$\begin{array}{cc}\mathrm{SO}(2m)/\mathrm{U}(m),\hfill & m5;\hfill \\ \mathrm{SU}(m)/\mathrm{SO}(m),\hfill & m3,m4;\hfill \\ \mathrm{SU}(2m)/\mathrm{Sp}(m),\hfill & m3;\hfill \\ \mathrm{Sp}(m)/\mathrm{U}(m),\hfill & m3.\hfill \end{array}$$
(Note that $`\mathrm{SO}(8)/\mathrm{U}(4)`$ is locally isometric to $`𝔾_2(^8)`$, $`\mathrm{SU}(4)/\mathrm{SO}(4)`$ is locally isometric to $`𝔾_3(^6)`$ and $`\mathrm{Sp}(2)/\mathrm{U}(2)`$ is locally isometric to $`𝔾_3(^5)`$, cf. , Ch. X, §6.4). In the sequel we will refer to these spaces as “structure spaces” since they can be interpreted as: a space of complex structures on $`^{2n}`$, spaces of real or quaternionic structures on $`^n`$ or $`^{2m}`$, respectively, and a space of complex structures on $`^n`$.
We will first consider the maximal subgroups $`H_1G`$; for classical groups $`G`$, they are given by Propositions 3.6, 3.7, and 3.8. Note that for $`\mathrm{SO}(n)`$ the subgroups (i) and (iii), for $`\mathrm{SU}(n)`$ the subgroups (i), (ii), (iii) and for $`\mathrm{Sp}(n)`$ the subgroups (i) and (ii) are symmetric, thus the actions of these groups are Hermann actions. The remaining types of subgroups are either given by tensor product representations or are simple irreducible subgroups. We will also study certain subactions of cohomogeneity one or transitive Hermann actions
Henceforth we will refer to the following maximal connected subgroups of the classical Lie groups (cf. Section 3) as (maximal) tensor product subgroups:
(8.2)
$$\begin{array}{cccc}H=\mathrm{SO}(p)\mathrm{SO}(q)\hfill & & G=\mathrm{SO}(pq),\hfill & p3,q3;\hfill \\ H=\mathrm{SU}(p)\mathrm{SU}(q)\hfill & & G=\mathrm{SU}(pq),\hfill & p3,q2;\hfill \\ H=\mathrm{SO}(p)\mathrm{Sp}(q)\hfill & & G=\mathrm{Sp}(pq),\hfill & p3,q1;\hfill \\ H=\mathrm{Sp}(p)\mathrm{Sp}(q)\hfill & & G=\mathrm{SO}(4pq),\hfill & p2,q1.\hfill \end{array}$$
###### Proposition 8.1 (Tensor product groups on “structure spaces”).
Let $`G`$ be a simple compact classical Lie group, let $`H_1`$ be a maximal tensor product subgroup of $`G`$ as in (8.2), let $`K`$ be a structure subgroup as in (8.1), and let $`HH_1`$ be a closed connected subgroup acting nontrivially on $`G/K`$. Then the $`H`$-action on $`G/K`$ is not polar.
###### Proof.
There are a few exceptions remaining not excluded by Proposition 5.6:
(8.3)
$$\begin{array}{ccc}\hfill H=\mathrm{SU}(3)\mathrm{SU}(2)& \text{acting on}& \mathrm{SU}(6)/\mathrm{Sp}(3)=G/K;\hfill \\ \hfill \mathrm{SU}(4)\mathrm{SU}(2)& \text{acting on}& \mathrm{SU}(8)/\mathrm{Sp}(4);\hfill \\ \hfill \mathrm{Sp}(3)\mathrm{Sp}(1)& \text{acting on}& \mathrm{SO}(12)/\mathrm{U}(6).\hfill \end{array}$$
We will apply Lemma 6.4 to show that none of the actions 8.3 can have a polar subaction. We use the information collected in Table 2 of to determine slice representations. In the case of the $`H=\mathrm{SU}(3)\mathrm{SU}(2)`$-action on $`G/K=\mathrm{SU}(6)/\mathrm{Sp}(3)`$, a slice representation is $`\mathrm{Ad}_{\mathrm{SO}(3)}\mathrm{Ad}_{\mathrm{SU}(2)}`$, which is polar. However, an explicit calculation shows that the normal space to a regular orbit is not a Lie triple system, thus by Proposition 4.1, the action is not polar. Let now $`UH`$ be closed connected proper subgroup acting polarly on $`G/K`$. By Theorem 6.1 and Lemma 2.6, the above slice representation restricted to $`(UK)_0`$ is either trivial, leading to a contradiction by Proposition 5.6 since the slice is $`9`$-dimensional, or equivalent to the action of $`HK`$, in which case $`U`$ must contain $`\mathrm{SO}(3)\mathrm{SU}(2)`$. But since $`\mathrm{SO}(3)\mathrm{SU}(3)`$ is maximal connected we have that the $`U`$-action on $`G/K`$ has a fixed point and is non-polar by Corollary 6.2.
Let us consider the $`\mathrm{SU}(4)\mathrm{SU}(2)`$-action on $`\mathrm{SU}(8)/\mathrm{Sp}(4)`$, a slice representation is $`\mathrm{Ad}_{\mathrm{SO}(4)}\mathrm{Ad}_{\mathrm{SU}(2)}`$, which is non-polar and polarity minimal, hence we may apply Lemma 6.4 (ii). For the $`\mathrm{Sp}(3)\mathrm{Sp}(1)`$-action we find a slice representation $`\mathrm{P}_2(\mathrm{Sp}(3))^2`$ of $`\mathrm{Sp}(3)\mathrm{U}(1)`$, which is also non-polar and polarity minimal, hence Lemma 6.4 (ii) also applies in this case. ∎
###### Proposition 8.2.
Let $`H_1G`$ and $`G/K`$ be as in Table 6 and let $`HH_1`$ be a closed connected subgroup acting nontrivially on $`G/K`$. Then the $`H`$-action on $`G/K`$ is not polar.
###### Proof.
None of the subgroups $`H_1G`$ fulfills the lower bound on its dimension given in Proposition 5.6, except the action of $`\mathrm{S}((\mathrm{U}(2)\mathrm{U}(3))\times \mathrm{U}(2))`$ on $`\mathrm{SU}(8)/\mathrm{Sp}(4)`$. However, an explicit calculation shows that the cohomogeneity of this action is $`12`$, hence this action is non-polar and polarity minimal by Lemma 3.3. ∎
###### Proposition 8.3 (Tensor product subgroups on Grassmannians).
Let $`G`$, $`H`$, $`K`$ be as in (8.4). Assume $`n=2,\mathrm{},\frac{pq}{2}`$ in cases (a), (b), (c). In case (d), assume that $`n=2,\mathrm{},2pq`$ and $`n3`$ if $`q=1`$. Assume further that in case (d) $`pq2`$. Then the action of $`H`$ on $`G/K`$ is non-polar and polarity minimal, i.e. any nontrivial action of a closed connected subgroup $`UH`$ on $`G/K`$ is non-polar.
(8.4)
$$\begin{array}{cccc}\text{(a)}& H=\mathrm{SO}(p)\mathrm{SO}(q),\hfill & G/K=\mathrm{SO}(pq)/\mathrm{SO}(n)\times \mathrm{SO}(pqn),\hfill & p3,q3;\hfill \\ \text{(b)}& H=\mathrm{SU}(p)\mathrm{SU}(q),\hfill & G/K=\mathrm{SU}(pq)/\mathrm{S}(\mathrm{U}(n)\times \mathrm{U}(pqn)),\hfill & p3,q2;\hfill \\ \text{(c)}& H=\mathrm{SO}(p)\mathrm{Sp}(q),\hfill & G/K=\mathrm{Sp}(pq)/\mathrm{Sp}(n)\times \mathrm{Sp}(pqn),\hfill & p3,q1;\hfill \\ \text{(d)}& H=\mathrm{Sp}(p)\mathrm{Sp}(q),\hfill & G/K=\mathrm{SO}(4pq)/\mathrm{SO}(n)\times \mathrm{SO}(4pqn),\hfill & p2,q1.\hfill \end{array}$$
###### Proof.
We use the slice representations which were explicitly determined in Section 2.3 of . In each case, one finds a non-polar, polarity minimal, and almost effective submodule of the slice representation, by Lemma 6.3 (iii) this implies that the slice representation is polarity minimal. But then Lemma 6.4 (iii) shows that the $`U`$-action is non-polar, since the normal space $`\mathrm{N}_{\mathrm{e}K}(H\mathrm{e}K)`$ contains a Lie triple system corresponding to an irreducible symmetric space of non-constant curvature in each case, as can be seen from the explicit description of the normal spaces in . ∎
###### Proposition 8.4.
Let $`H_1G`$ and $`G/K`$ be as in Table 7 and let $`HH_1`$ be a closed connected subgroup acting nontrivially on $`G/K`$. Then the $`H`$-action on $`G/K`$ is not polar.
###### Proof.
Consider the first four items in Table 7. By Lemma 11.2, the corresponding Hermann action (indicated in the first column) has a totally geodesic orbit isometric to $`𝔾_k(^{pq})`$, $`𝔾_k(^{pq})`$, $`𝔾_k(^{4pq})`$, and $`𝔾_k(^{pq})`$, respectively, on which $`H`$ acts. It follows from Proposition 8.3 and Lemma 4.2 that $`H`$ acts non-polarly except if $`p=3`$, $`q=1`$, $`k=2`$ in case of the fourth action; however in this case the normal space of a principal orbits is not a Lie triple system and $`H_1`$ acts non-polarly by Proposition 4.1; closed proper subgroups of $`H_1`$ are excluded by Proposition 5.6 (iii).
Let us now consider the last two items of Table 7. Assume $`H`$ acts polarly and nontrivially on $`G/K`$. Then the $`H`$-action has a singular orbit by Corollary 5.3. As can be seen from Table 10, the almost effective slice representation of the $`H_1`$-action on $`G/K`$ occurs also as a submodule of the isotropy representation $`\chi (H_1,H_1K)`$ of the $`H_1`$-orbit $`H_1\mathrm{e}K`$. Thus the $`H`$-action on $`G/K`$ is non-polar by Proposition 6.3 (ii). ∎
## 9. Subactions of simple irreducible groups
We will now study simple irreducible maximal subgroups of the classical groups acting on the classical symmetric spaces. We start with actions on the spaces (8.1).
###### Proposition 9.1 (Simple irreducible groups on “structure spaces”).
Let $`G`$ be a simple compact classical Lie group $`\mathrm{SO}(n)`$, $`\mathrm{SU}(n)`$ or $`\mathrm{Sp}(n)`$ and let $`\rho :HG`$ be an irreducible representation of corresponding (real, complex or quaternionic) type where $`H`$ is a simple compact Lie group and such that $`\rho (H)`$ is a maximal connected subgroup of $`G`$. Let $`KG`$ be a subgroup as in (8.1) such that $`\mathrm{rk}(G/K)2`$. Then the action of any closed subgroup of $`\rho (H)`$ on $`G/K`$ is non-polar except for the Hermann actions of subgroups conjugate to $`\mathrm{Spin}(7)\mathrm{SO}(8)`$ on $`\mathrm{SO}(8)/\mathrm{U}(4)`$.
###### Proof.
For the spaces $`\mathrm{SU}(n)/\mathrm{SO}(n)`$ and $`\mathrm{Sp}(n)/\mathrm{U}(n)`$ this follows directly from Lemmata 2.7, 2.8 of and Proposition 5.6.
Let us consider the spaces $`\mathrm{SO}(n)/\mathrm{U}(\frac{n}{2})`$, $`n8`$. By Proposition 5.6 we have that $`dim(H)\frac{n^2}{4}n`$ if $`H`$ acts polarly on $`\mathrm{SO}(n)/\mathrm{U}(\frac{n}{2})`$. Hence $`\rho `$ is a representation as described in Lemma 2.6 (i), (iv) of and all possibilities for $`\rho `$ are given in the table of , Lemma 2.8 (i). However, all of these subgroups $`\rho (H)\mathrm{SO}(n)`$ are excluded by the dimension bounds given in Proposition 5.6, except $`\mathrm{Spin}(7)\mathrm{SO}(8)`$.
For the spaces $`\mathrm{SU}(n)/\mathrm{Sp}(\frac{n}{2})`$, $`n3`$, it follows from Proposition 5.6 that $`dim(H)\frac{n^2}{2}2n`$ for a group $`H`$ acting polarly on $`\mathrm{SU}(n)/\mathrm{Sp}(\frac{n}{2})`$. Thus $`\rho `$ is a representation as in Lemma 2.6 (ii), (iv) of and all such representations $`\rho `$ are given by the table in , Lemma 2.8 (ii). However, none of the simple groups there fulfills the necessary condition on its dimension given by Proposition 5.6. ∎
We also need to consider certain subactions of cohomogeneity one or transitive actions.
###### Lemma 9.2.
1. Let $`k\{1,2,3\}`$. Let $`H`$ be a simple compact connected Lie group and let $`\rho :H\mathrm{SO}(2nk)`$ be an irreducible representation of real type such that $`\rho (H)\mathrm{SO}(2nk)`$ is maximal connected. Then any closed subgroup of $`\rho (H)\times \mathrm{SO}(k)`$ acts non-polarly on $`\mathrm{SO}(2n)/\mathrm{U}(n)`$, except if $`\rho `$ is equivalent to the $`7`$-dimensional irreducible representation of $`\mathrm{G}_2`$ and $`k=1`$.
2. Let $`k\{1,2,3\}`$. Let $`H`$ be a simple compact connected Lie group and let $`\rho :H\mathrm{SU}(2nk)`$ be an irreducible representation of complex type such that $`\rho (H)\mathrm{SU}(2nk)`$ is maximal connected. Then any closed subgroup of $`\mathrm{S}((\rho (H)\mathrm{U}(1))\times \mathrm{U}(k))`$ acts non-polarly on $`\mathrm{SU}(2n)/\mathrm{Sp}(n)`$.
3. Let $`H`$ be a simple compact connected Lie group and let $`\rho :H\mathrm{SU}(n)`$ be an irreducible representation of complex type such that $`\rho (H)\mathrm{SU}(n)`$ is maximal connected. Then any closed subgroup of $`\mathrm{S}((\rho (H)\mathrm{U}(1))\times \mathrm{U}(1))`$ acts non-polarly on $`\mathrm{SU}(n)/\mathrm{SO}(n)`$.
4. Let $`H`$ be a simple compact connected Lie group and let $`\rho :H\mathrm{Sp}(n)`$ be an irreducible representation of real type such that $`\rho (H)\mathrm{Sp}(n)`$ is maximal connected. Then any closed subgroup of $`\rho (H)\times \mathrm{Sp}(1)`$ acts non-polarly on $`\mathrm{Sp}(n)/\mathrm{U}(n)`$.
###### Proof.
The proof is almost literally the same as the proof of Proposition 9.1. For the cases (i) and (ii), we may use the tables in parts (i) and (ii) of Lemma 2.8 in . The only representation not excluded by this argument are the $`8`$-dimensional spin representation of $`\mathrm{Spin}(7)`$ and the $`7`$-dimensional representation of $`\mathrm{G}_2`$. However, in case of the actions of $`\mathrm{Spin}(7)\times \mathrm{SO}(2)`$ and $`\mathrm{G}_2\times \mathrm{SO}(3)`$ on $`\mathrm{SO}(10)/\mathrm{U}(5)`$, the normal space at a principal orbit is not a Lie triple system and hence these actions are non-polar by Proposition 4.1. Closed connected subgroups of these groups can be shown to act non-polarly by the same argument or are excluded by Proposition 5.6. The action of $`\mathrm{G}_2`$ on $`\mathrm{SO}(8)/\mathrm{U}(4)`$ is orbit equivalent to the action of $`\mathrm{Spin}(7)`$ on $`\mathrm{SO}(8)/\mathrm{U}(4)`$. The statements (iii) and (iv) follow directly from Lemmata 2.7 and 2.8 of . ∎
We will now consider the maximal simple irreducible subgroups of the classical groups $`\mathrm{SO}(n)`$, $`\mathrm{SU}(n)`$, $`\mathrm{Sp}(n)`$, given by irreducible representations of the real, complex, or quaternionic type, respectively, and their actions on the corresponding Grassmannians $`𝔾_k(𝕂^n)`$, $`𝕂=`$, $``$ or $``$, of higher rank. A necessary condition for polarity on the dimension of these subgroups is given by Propositions 5.6 and 5.5. The irreducible representations of simple compact Lie groups whose degrees are sufficiently low can be obtained from Lemma 2.6 of , see also the tables in Lemma 2.8 and the Appendix of . These representations are given by Table 8; the column marked with $`k_{\mathrm{max}}`$ indicates the maximal rank $`k\frac{n}{2}`$ for which the necessary condition for polarity of an action on $`𝔾_k(𝕂^n)`$ given by Proposition 5.6 is fulfilled. We only list such representations of complex or quaternionic type where $`k_{\mathrm{max}}2`$ and representations of real type where $`k_{\mathrm{max}}3`$, since polar actions on $`𝔾_2(^n)`$ have been classified in . It turns out there are no such representations of complex or quaternionic type.
#### Subactions of $`\mathrm{Ad}(\mathrm{SU}(3))`$ on $`𝔾_3(^8)`$ and $`𝔾_4(^8)`$
The group $`H=\mathrm{SU}(3)`$ acts on its Lie algebra h by the adjoint representation $`\mathrm{Ad}:H\mathrm{SO}(\text{h})`$ and we obtain a subgroup $`\mathrm{Ad}(\mathrm{SU}(3))\mathrm{SO}(8)`$ by identifying h with $`^8`$. We will study the actions of this group on the Grassmannians $`𝔾_k(^8)`$, $`k=3,4`$. Any closed connected proper subgroups of $`\mathrm{SU}(3)`$ are of dimension $`4`$ by Table 3 of and are thus excluded by Proposition 5.6.
The maximal connected subgroup $`\mathrm{SO}(3)\mathrm{SU}(3)`$ leaves a $`3`$-dimensional subspace of $`^8`$ invariant, thus it is the connected component of an isotropy subgroup of the $`\mathrm{Ad}(\mathrm{SU}(3))`$-action on $`𝔾_3(^8)`$. The slice representation contains the irreducible $`7`$-dimensional representation of $`\mathrm{SO}(3)`$ and is hence non-polar .
Consider now the subgroup $`\mathrm{S}(\mathrm{U}(1)\times \mathrm{U}(2))\mathrm{SU}(3)`$. The action of $`H`$ on $`\text{h}=^8`$ restricted to $`\mathrm{S}(\mathrm{U}(1)\times \mathrm{U}(2))`$ leaves the $`4`$-dimensional linear subspace corresponding to $`𝔰(𝔲(1)+𝔲(2))\text{h}`$ invariant. Thus the maximal connected $`\mathrm{S}(\mathrm{U}(1)\times \mathrm{U}(2))`$ coincides with the connected component of the stabilizer of the $`H`$-action on $`𝔾_4(^8)`$. Its slice representation contains an $`8`$-dimensional irreducible representation of $`\mathrm{S}(\mathrm{U}(1)\times \mathrm{U}(2))`$ which is non polar and the $`H`$-action on $`𝔾_4(^8)`$ is polarity minimal by Lemma 6.4 (ii).
#### Subactions of $`\mathrm{Spin}(7)`$ on $`𝔾_3(^8)`$ and $`𝔾_4(^8)`$
The subgroup $`\mathrm{Spin}(7)\mathrm{SO}(8)`$ gives rise to a Hermann action since its Lie algebra is the fixed point set of an involution of $`𝔰𝔬(8)`$.
#### Subactions of $`\mathrm{Spin}(9)`$ on $`𝔾_3(^{16})`$
The action on $`𝔾_3(^{16})`$ was shown not to be polar in , the slice representation being equivalent to a $`16`$-dimensional non-polar irreducible representation of $`\mathrm{Sp}(1)\mathrm{Sp}(2)`$. Thus by Lemma 6.4 (ii), no subaction of the $`\mathrm{Spin}(9)`$-action on $`𝔾_3(^{16})`$ is polar.
#### Subactions of $`\mathrm{Spin}(9)`$ on $`𝔾_4(^{16})`$
Consider the subgroup $`H_0=\mathrm{Spin}(4)\mathrm{Spin}(4)\mathrm{Spin}(8)\mathrm{Spin}(9)`$, its action on $`^{16}`$ leaves a four-dimensional subspace $`V`$ invariant. Since $`H_0\mathrm{Spin}(9)`$ is a subgroup of maximal rank, it is easy to check that no other connected subgroup of $`\mathrm{Spin}(9)`$ containing $`H_0`$ leaves $`V`$ invariant and thus $`H_0`$ is the connected component of an isotropy subgroup of the $`\mathrm{Spin}(9)`$-action on $`^{16}`$. The slice representation is equivalent to the sum of two $`12`$-dimensional irreducible modules and is easily seen to be non-polar and polarity minimal.
#### Subactions of $`\mathrm{Sp}(3)`$ on $`𝔾_3(^{14})`$
An isotropy subgroup of the $`\mathrm{Sp}(3)`$-action on $`𝔾_3(^{14})`$ is
$$\mathrm{SO}(3)\mathrm{U}(1)\mathrm{U}(3)\mathrm{Sp}(3).$$
Its $`16`$-dimensional slice representation does not contain any trivial submodule and is therefore non-polar by . Any proper subgroups of $`\mathrm{Sp}(3)`$ can be excluded by Proposition 5.6.
#### Subactions of $`\mathrm{F}_4`$ on $`𝔾_3(^{26})`$
The maximal connected subgroups of maximal rank in $`H_1=\mathrm{F}_4`$ are, see , Chapter 1, § 3.11.
(9.1)
$$\mathrm{Sp}(3)\mathrm{Sp}(1),\mathrm{SU}(3)\mathrm{SU}(3),\mathrm{Spin}(9)$$
We will determine a slice representation for the $`H_1`$-action on $`𝔾_3(^{26})`$. According to , Table 25, p. 199, the subgroup $`\mathrm{Spin}(9)`$ acts on $`^{26}`$ by the direct sum of the $`9`$-dimensional standard representation, the $`16`$-dimensional spin representation, and a one-dimensional trivial representation. Thus if we further restrict this representation to the maximal connected subgroup $`\mathrm{Spin}(7)\mathrm{SO}(2)`$ of $`\mathrm{Spin}(9)`$, a three-dimensional subspace $`W`$ is left invariant and it follows that $`\mathrm{Spin}(7)\mathrm{SO}(2)`$ is contained in an isotropy subgroup $`(H_1)_W`$ of the $`H_1`$-action on $`𝔾_3(^{26})`$. The subgroup $`\mathrm{Spin}(7)\mathrm{SO}(2)\mathrm{F}_4`$ is of maximal rank and it can be deduced from Table 25 of that none of the groups in (9.1) leaves a three-dimensional subspace of $`^{26}`$ invariant. Hence $`\mathrm{Spin}(7)\mathrm{SO}(2)`$ is the connected component of the isotropy subgroup $`(H_1)_W`$ and the slice representation is, by a dimension count, equivalent to $`^72^2^8`$, where $`\mathrm{Spin}(7)`$ acts on $`^8`$ by the spin representation, hence the action of $`\mathrm{F}_4`$ on $`𝔾_3(^{26})`$ is non-polar and polarity minimal by Proposition 6.3 (ii) and Lemma 6.4 (ii).
#### Subactions of $`\mathrm{G}_2`$ on $`𝔾_3(^7)`$
This is a cohomogeneity one action and its subactions will be treated in Section 12.
## 10. Polar actions on the exceptional spaces
In this section we will study those isometric actions on the exceptional symmetric spaces of compact type which are subactions neither of Hermann actions nor of cohomogeneity one actions. It will turn out that none of these actions is polar.
The maximal connected subgroups of the simple compact Lie groups were determined in , Tables 12 and 12a, p. 150–151, and Theorem 14.1, p. 231. By Theorem 5.4 and Lemma 3.3, the cohomogeneity of a polar action on a symmetric quotient $`G/K`$ of a simple Lie group $`G`$ is at most $`\mathrm{rk}(G)+\mathrm{rk}(K)`$. By Proposition 5.5, this estimate can be further improved for Hermitian symmetric spaces, for which the cohomogeneity is at most $`\mathrm{rk}(G)`$. From this it follows by using the classification of symmetric spaces, see Table 4, that a group acting polarly on a symmetric quotient $`G/K`$ with $`\mathrm{rk}(G/K)2`$ of one of the simple exceptional Lie groups $`G=\mathrm{E}_6`$, $`\mathrm{E}_7`$, $`\mathrm{E}_8`$, $`\mathrm{F}_4`$, $`\mathrm{G}_2`$ is at least of dimension $`16`$, $`47`$, $`96`$, $`20`$, $`4`$, respectively. (We do need not consider the Cayley plane $`\mathrm{F}_4/\mathrm{Spin}(9)`$, since it is of rank one.) First we would like to recall a method to describe certain subgroups of a (semi)simple compact Lie group in terms of the root system, which is particularly useful for our purposes, see , § 8.3 and , Ch. 1, § 3.11.
###### Remark 10.1 (Borel-De Siebenthal theory).
Let $`G`$ be a connected compact simple Lie group. A subgroup $`HG`$ is called a subgroup of maximal rank if $`\mathrm{rk}(H)=\mathrm{rk}(G)`$, i.e. $`H`$ contains a maximal torus $`T`$ of $`G`$. Consider the root space decomposition $`\text{g}_{}=\text{g}_0+_{\alpha R}\text{g}_\alpha `$, where $`\text{g}_0`$ is the complexification of the maximal abelian subalgebra of g tangent to $`T`$. Since the Lie algebra $`\text{h}_{}`$ contains $`\text{g}_0`$, it is a $`\text{g}_0`$-stable subspace of $`\text{g}_{}`$, and it follows that $`\text{h}_{}=\text{g}_0+_{\alpha S}\text{g}_\alpha `$ where $`SR`$ is a subset of the root system. Conversely, from suitable subsets $`SR`$, one may construct the Lie algebra of a subgroup $`HG`$ of maximal rank, see , Chapter 1, § 3.11. In particular, one can obtain all maximal connected subgroups of maximal rank by such a construction. These are obtained by deleting certain vertices from the extended Dynkin diagram, see for details. The classification of all such subgroups up to conjugation by automorphisms of $`G`$ is given in Table 5, p. 64 of or in Table 12, p. 150 of , see also .
Now assume $`H`$ and $`K`$ are both subgroups of maximal rank in $`G`$. Then we can use the above description to obtain information about the $`H`$-action on the homogeneous space $`G/K`$, in particular, to compute an isotropy algebra together with its slice representation. In fact, we may assume by conjugation of $`K`$ with a suitable element from $`G`$ that both $`H`$ and $`K`$ contain a maximal torus $`T`$ of $`G`$. Then $`\text{h}_{}=\text{g}_0+_{\alpha S}\text{g}_\alpha `$ and $`\text{k}_{}=\text{g}_0+_{\alpha S^{}}\text{g}_\alpha `$ for some subsets $`S,S^{}`$ of the root system $`R`$. In particular, the complexified isotropy algebra $`(\text{h}\text{k})_{}`$ of the $`H`$-action on $`G/K`$ at $`\mathrm{e}K`$ is spanned by the Cartan algebra $`\text{g}_0`$ and the root spaces corresponding to the roots in the intersection $`SS^{}`$ and it follows that $`HK`$ is also a subgroup of maximal rank in $`G`$. On the other hand, the complexified normal space $`(\text{h}^{}\text{k}^{})_{}`$ of the $`H`$-orbit through $`\mathrm{e}K`$ is spanned by the root spaces corresponding to the roots in $`R\left(SS^{}\right)`$. Since $`T`$ is also a maximal torus of $`HK`$, the roots in $`R\left(SS^{}\right)`$ are exactly the weights of the slice representation of $`HK`$ on the normal space $`\text{h}^{}\text{k}^{}`$. It follows that the $`H`$-orbit through $`\mathrm{e}K`$ is a singular orbit, since $`T`$ acts nontrivially on $`\text{h}^{}\text{k}^{}`$, in fact, the slice representation does not have any trivial submodules, since the complexified normal space is spanned by root spaces corresponding to non-zero roots. In the special case $`H=K`$, one obtains the isotropy representation $`\chi (G,K)`$ by this method.
Note that if a subgroup $`HG`$ is a fixed point set of an inner automorphism $`\sigma `$ of $`G`$, i.e. $`\sigma (x)=gxg^1`$, it is a subgroup of maximal rank, since the element $`g=\mathrm{exp}(X)`$, $`X\text{g}`$, lies in a maximal torus $`T`$ of $`G`$ and it follows that $`\text{h}_{}=\text{g}_0+_{\{\alpha X\mathrm{ker}\alpha \}}\text{g}_\alpha `$ where $`\text{g}_0`$ is the complexified Lie algebra of $`T`$. (Conversely, if a subgroup of maximal rank is the fixed point set of an automorphism, then the automorphism is inner.)
Let $`H`$ and $`K`$ be two subgroups of maximal rank with common maximal torus $`T`$ as above. If both groups are fixed point sets of involutions i.e. $`H=G^\sigma `$, $`K=G^\tau `$, then it follows that the involutions $`\sigma `$ and $`\tau `$ commute, since they both act as either plus or minus identity on the root spaces of g. This shows that if $`\sigma `$ and $`\tau `$ are two inner involutions of a simple compact Lie group $`G`$, then $`\tau `$ is conjugate to an involution which commutes with $`\sigma `$, cf. .
### 10.1. Symmetric spaces of $`\mathrm{E}_6`$
The maximal connected non-symmetric subgroups of $`\mathrm{E}_6`$ of dimension $`16`$ are $`\mathrm{SU}(3)\mathrm{SU}(3)\mathrm{SU}(3)`$ and $`\mathrm{G}_2^1\mathrm{A}_{2}^{2}{}_{}{}^{\prime \prime }`$, see . (The upper indices denote the Dynkin index of subgroups and the primes are used to distinguish non-conjugate subgroups of the same Dynkin index). By a dimension count, no closed subgroup of these groups acts polarly on the spaces $`\mathrm{E}_6/(\mathrm{Sp}(4)/\{\pm 1\})`$, $`\mathrm{E}_6/\mathrm{SU}(6)\mathrm{Sp}(1)`$ or $`\mathrm{E}_6/\mathrm{Spin}(10)\mathrm{U}(1)`$.
It remains to determine the polar actions on $`\mathrm{E}_6/\mathrm{F}_4`$. We start with the group $`H_1=\mathrm{SU}(3)\mathrm{SU}(3)\mathrm{SU}(3)`$. The subgroup $`\mathrm{SU}(3)\mathrm{SU}(3)\mathrm{SU}(3)`$ is constructed from the extended Dynkin diagram of $`\mathrm{E}_6`$ as follows, cf. Remark 10.1.
The vertices numbered $`1,\mathrm{},6`$ correspond to the simple roots $`\alpha _1,\mathrm{},\alpha _6`$ of $`\mathrm{E}_6`$ and the vertex with number $`0`$ represents $`\alpha _0`$, where $`\alpha _0`$ is the maximal root. Now the group $`H_1=\mathrm{SU}(3)\mathrm{SU}(3)\mathrm{SU}(3)`$ arises from the extended Dynkin diagram if one deletes the central vertex $`3`$, i.e. it is the regular subgroup whose simple roots are $`\alpha _1`$, $`\alpha _2`$, $`\alpha _4`$, $`\alpha _5`$, $`\alpha _6`$, $`\alpha _0`$. The subgroup $`\mathrm{F}_4\mathrm{E}_6`$ is the fixed point set of the diagram automorphism $`\sigma `$ of $`\mathrm{E}_6`$ which maps $`\alpha _1\alpha _5`$, $`\alpha _2\alpha _4`$, $`\alpha _4\alpha _2`$, $`\alpha _5\alpha _1`$ and leaves $`\alpha _3`$ and $`\alpha _6`$ fixed. This automorphism $`\sigma `$ also leaves $`\alpha _0`$ fixed, since $`\alpha _0=\alpha _1+2\alpha _2+3\alpha _3+2\alpha _4+\alpha _5+2\alpha _6`$, see , Chapter 1, § 3.11. It follows that $`\sigma `$ also acts on $`H_1`$, i.e. trivially on one $`\mathrm{SU}(3)`$-factor (the one whose simple roots are $`\alpha _6`$, $`\alpha _0`$) and by interchanging the other two $`\mathrm{SU}(3)`$-factors. Thus $`H_1\mathrm{F}_4`$ is the fixed point set $`H_1^\sigma `$ and is hence isomorphic to $`\mathrm{SU}(3)\mathrm{\Delta }\mathrm{SU}(3)`$, where the $`\mathrm{\Delta }\mathrm{SU}(3)`$-factor is diagonally embedded into two of the $`\mathrm{SU}(3)`$-factors of $`H_1`$. Let us determine the slice representation of the $`H_1`$-action on $`M`$, it is a submodule of $`\chi (\mathrm{E}_6,H_1)`$ restricted to $`H_1\mathrm{F}_4`$. The real $`54`$-dimensional isotropy representation $`\chi (\mathrm{E}_6,H_1)`$ is, after complexification,
(10.1)
$$\left(\text{}\right)\left(\text{}\right),$$
see , Corollary 13.2, i.e. the isotropy representation is equivalent to the action of $`\mathrm{SU}(3)\mathrm{SU}(3)\mathrm{SU}(3)`$ on $`^3^3^3`$ by the tensor product of the standard representations. If we restrict this representation to the subgroup $`\mathrm{SU}(3)\mathrm{\Delta }\mathrm{SU}(3)`$, it splits into the irreducible modules $`(^3\mathrm{Sym}^2^3)`$ and $`(^3\mathrm{\Lambda }^2^3)`$, where the first $`\mathrm{SU}(3)`$-factor acts on $`^3`$ and the $`\mathrm{\Delta }\mathrm{SU}(3)`$-factor acts on $`\mathrm{Sym}^2^3`$ or $`\mathrm{\Lambda }^2^3`$, respectively. A dimension count shows that the real $`36`$-dimensional slice representation of the $`H_1`$-action is equivalent to the first irreducible summand, hence it is non-polar and polarity minimal. We conclude that the $`H_1`$-action on $`G/K`$ is non-polar and polarity minimal by Lemma 6.4 (ii).
Now consider subactions of $`H_1=\mathrm{G}_2^1\mathrm{A}_{2}^{2}{}_{}{}^{\prime \prime }`$ on $`G/K=\mathrm{E}_6/\mathrm{F}_4`$. We determine a slice representation of the $`H_1`$-action on $`G/K`$. First observe that $`\mathrm{F}_4`$ contains the subgroup $`\mathrm{G}_2^1`$, according to Table 39 of , p. 233. Since the subgroup $`\mathrm{F}_4\mathrm{E}_6`$ has Dynkin index $`1`$, it follows that $`\mathrm{G}_2^1\mathrm{E}_6`$ also has Dynkin index $`1`$, see , Ch. I, § 2. By Table 25 of , p. 200, there is only one conjugacy class of subgroups isomorphic to $`\mathrm{G}_2`$ of Dynkin index $`1`$ in $`\mathrm{E}_6`$ and it follows that an isotropy subgroup $`(H_1)_x`$ of the $`H_1`$-action on $`G/K`$ contains $`\mathrm{G}_2^1`$. The homogeneous space $`G/H_1`$ is strongly isotropy irreducible, see , Theorem 3.1, p. 66, and its isotropy representation decomposes into $`8`$ equivalent $`7`$-dimensional irreducible modules when restricted to $`\mathrm{G}_2`$. Thus the dimension of the normal space $`\mathrm{N}_x(H_1x)`$ to the $`H_1`$-orbit through $`x`$ is a multiple of $`7`$. The only possibility is a $`21`$-dimensional slice representation which splits into $`3`$ irreducible $`7`$-dimensional modules when restricted to $`\mathrm{G}_2^1`$. By and , such a representation is non-polar and it is polarity minimal by Lemma 6.3, part (i) or (ii), and hence we can apply Lemma 6.4 to show that no closed subgroup $`HH_1`$ acts polarly on $`G/K`$.
### 10.2. Symmetric spaces of $`\mathrm{E}_7`$
The only maximal connected non-symmetric subgroup of dimension $`47`$ is $`H_1=\mathrm{F}_4^1\mathrm{A}_{1}^{3}{}_{}{}^{\prime \prime }`$, see . By Lemma 3.3 and Proposition 5.5, respectively, no closed subgroup of $`H_1`$ acts polarly on $`\mathrm{E}_7/(\mathrm{SU}(8)/\{\pm 1\})`$ or on $`\mathrm{E}_7/\mathrm{E}_6\mathrm{U}(1)`$.
Let us determine an isotropy subgroup of the $`H_1`$-action on $`G/K=\mathrm{E}_7/\mathrm{SO}^{}(12)\mathrm{Sp}(1)`$. First observe that $`\text{h}_1`$ contains a subalgebra $`𝔰𝔭𝔦𝔫(9)𝔣_4\text{h}_1`$. By Table 25 of , p. 201, there is only one conjugacy class of subalgebras isomorphic to $`𝔰𝔭𝔦𝔫(9)`$ in $`𝔢_7`$ and it follows that, after conjugation, this subalgebra coincides with the subalgebra $`𝔰𝔭𝔦𝔫(9)𝔰𝔭𝔦𝔫(12)K`$. Thus there is an isotropy subgroup $`(H_1)_x`$ of the $`H_1`$-action on $`G/K`$ whose Lie algebra contains $`𝔰𝔭𝔦𝔫(9)`$ as a subalgebra. The $`64`$-dimensional isotropy representation of $`\mathrm{E}_7/\mathrm{SO}^{}(12)\mathrm{Sp}(1)`$ decomposes into $`4`$ copies of the $`16`$-dimensional spin representation when restricted to $`\mathrm{Spin}(9)`$. Thus the dimension of an orbit $`H_1x`$ is a multiple of $`16`$ and it follows by a dimension count that the Lie algebra of $`(H_1)_x`$ is isomorphic to $`𝔰𝔭𝔦𝔫(9)𝔞_1`$. The $`32`$-dimensional slice representation at the point $`x`$ is the sum of two modules equivalent to the $`16`$-dimensional spin representation of $`\mathrm{Spin}(9)`$, hence non-polar and polarity minimal and thus by Lemma 6.4, the $`H_1`$-action is non-polar and polarity minimal.
Now consider the action of $`H_1=\mathrm{F}_4^1\mathrm{A}_{1}^{3}{}_{}{}^{\prime \prime }`$ on $`G/K=\mathrm{E}_7/\mathrm{E}_6\mathrm{U}(1)`$. By Table 25 of , p. 201, there is only one conjugacy class of subgroups isomorphic to $`\mathrm{F}_4`$ in $`\mathrm{E}_7`$ and it follows that an isotropy subgroup $`(H_1)_x`$ of the $`H_1`$-action on $`G/K`$ contains $`\mathrm{F}_4`$. The space $`G/H_1`$ is strongly isotropy irreducible, see , Theorem 3.1, and its isotropy representation decomposes into $`3`$ equivalent $`26`$-dimensional irreducible modules when restricted to $`\mathrm{F}_4`$. This shows that the dimension of the normal space $`\mathrm{N}_x(H_1x)`$ is a multiple of $`26`$ and it follows by a dimension count that the connected component of the isotropy subgroup $`(H_1)_x`$ is isomorphic to $`\mathrm{F}_4\mathrm{U}(1)`$. Thus the slice representation is non-polar and polarity minimal and we can apply Lemma 6.4 (i) to show that no closed subgroup $`HH_1`$ acts polarly on $`G/K`$.
### 10.3. Symmetric spaces of $`\mathrm{E}_8`$
The maximal connected subgroups of $`\mathrm{E}_8`$ whose dimension is at least $`96`$ are symmetric. Hence any polar action on the symmetric spaces $`\mathrm{E}_8/\mathrm{SO}^{}(16)`$ and $`\mathrm{E}_8/\mathrm{E}_7\mathrm{Sp}(1)`$ are subactions of Hermann actions.
### 10.4. Actions on $`\mathrm{F}_4/\mathrm{Sp}(3)\mathrm{Sp}(1)`$
All maximal connected subgroups of $`\mathrm{F}_4`$ whose dimension is at least $`20`$ are symmetric. Thus any polar action on the space $`\mathrm{F}_4/\mathrm{Sp}(3)\mathrm{Sp}(1)`$ is a subaction of a Hermann action.
### 10.5. Actions on $`\mathrm{G}_2/\mathrm{SO}(4)`$
The maximal connected subgroups of $`\mathrm{G}_2`$ are
(10.2)
$$\mathrm{SO}(4),\mathrm{SU}(3),\mathrm{A}_1^{28},$$
where $`\mathrm{A}_1^{28}`$ is a maximal connected subgroup in $`\mathrm{G}_2`$ of type $`\mathrm{A}_1`$, cf. . If the group $`H`$ acting on $`M=\mathrm{G}_2/\mathrm{SO}(4)`$ is contained in $`\mathrm{SO}(4)`$, then the action has a fixed point. The group $`\mathrm{SU}(3)`$ acts with cohomogeneity one on $`M`$. The only closed connected subgroup $`H\mathrm{SU}(3)`$ of dimension $`4`$ is $`\mathrm{S}(\mathrm{U}(1)\times \mathrm{U}(2))\mathrm{U}(2)`$, whose action on $`M`$ has a fixed point. Subgroups of rank one are ruled out by a dimension count.
## 11. Subactions of hyperpolar actions
To complete the classification, it remains to study subactions of cohomogeneity one and transitive actions. We will need the following lemmata to study subactions of Hermann actions. The first lemma shows that the slice representations of a Hermann action are s-representations.
###### Lemma 11.1 (Slice representations of Hermann actions).
Let $`G`$ be a connected simple compact Lie group and let $`\sigma `$, $`\tau `$ be two involutive automorphisms of $`G`$. Let $`K=G_0^\sigma `$ be the connected component of the fixed point set of $`\sigma `$ and let $`H_1=G_0^\tau `$ be the connected component of the fixed point set of $`\tau `$. Consider the $`H_1`$-action on $`G/K`$. Then the exponential image $`S=\mathrm{exp}_{\mathrm{e}K}(\mathrm{N}_{\mathrm{e}K}(H_1\mathrm{e}K))`$ of the normal space to the orbit through $`\mathrm{e}K`$ is a totally geodesic submanifold locally isometric to a symmetric space $`G^{\sigma \tau }/G^\sigma G^\tau `$ whose isotropy representation is on the Lie algebra level equivalent to the slice representation of the $`H_1`$-action on $`G/K`$ at $`\mathrm{e}K`$.
###### Proof.
The Lie algebra g of $`G`$ admits the two decompositions
(11.1)
$$\text{g}=\text{k}\text{p}=\text{h}_1\text{m}_1,$$
where p and $`\text{m}_1`$ are the $`1`$-eigenspaces of $`\sigma _{}`$ and $`\tau _{}`$, respectively. The normal space $`\mathrm{N}_{\mathrm{e}K}(H_1\mathrm{e}K)=\text{p}\text{m}_1\text{p}`$ is a Lie triple system and $`\text{k}\text{h}_1\text{p}\text{m}_1`$ is the Lie algebra generated by $`\text{p}\text{m}_1`$; the isotropy algebra at $`\mathrm{e}K`$ of the $`H_1`$-action on $`G/K`$ is just $`\text{h}_1\text{k}`$ and its action on the normal space $`\text{p}\text{m}_1`$ agrees on the Lie algebra level with the isotropy representation of $`G^{\sigma \tau }/G^\sigma G^\tau `$, which is locally isometric to $`S`$ by Proposition 3.2. ∎
In the special case where the two involutions defining a Hermann action commute (possibly after conjugation), the action has a totally geodesic orbit. The pairs of involutions on the compact simple Lie groups for which this is the case have been determined in . By Lemma 4.2, a polar subaction acts also polarly on this totally geodesic orbit.
###### Lemma 11.2 (Subactions of Hermann actions with commuting involutions).
Let $`G`$, $`K`$, and $`H_1`$ be as in Lemma 11.1. Assume in addition that $`\sigma \tau =\tau \sigma `$. Let $`HH_1`$ be a closed connected subgroup acting polarly on $`M=G/K`$. Then the $`H_1`$-orbit $`H_1\mathrm{e}K=H_1/H_1K`$ through $`\mathrm{e}KM`$ is a totally geodesic submanifold and $`H`$ acts polarly on the symmetric space $`H_1\mathrm{e}KH_1/H_1K`$.
###### Proof.
Since $`\sigma `$ and $`\tau `$ commute, we have the direct sum decomposition
(11.2)
$$\text{g}=(\text{k}\text{h}_1)(\text{k}\text{m}_1)(\text{p}\text{h}_1)(\text{p}\text{m}_1).$$
Consider now the $`H_1`$-action on the symmetric space $`G/K`$. We can identify p with the tangent space $`\mathrm{T}_{\mathrm{e}K}G/K`$. Then $`\text{h}_1\text{p}`$ is the tangent space of the $`H_1`$-orbit through the point $`\mathrm{e}K`$. Using the Cartan relations for the decompositions (11.1), it is easy to verify that $`\text{h}_1\text{p}`$ is a Lie triple system. Hence the $`H_1`$-orbit through $`\mathrm{e}K`$ is totally geodesic by Proposition 3.2. Clearly, the action of $`H`$ leaves all $`H_1`$-orbits invariant and the polarity of the $`H`$-action on $`H_1\mathrm{e}K`$ follows from Lemma 4.2. ∎
The following lemma is a just simple reformulation of the criterion for polarity given by Proposition 4.1 in the special case of subaction of a Hermann action; it is, however, useful in particular to study polar actions on the exceptional spaces since it enables us to test for polarity on a subspace.
###### Lemma 11.3.
Let $`G`$, $`K`$, $`H_1`$, $`H`$ and $`M=G/K`$ be as in Lemma 11.1. Assume the group $`H`$ acts transitively on the $`H_1`$-orbit through $`\mathrm{e}K`$. Then the action of $`H`$ on $`G/K`$ is polar if and only if the action of the action of $`HK`$ on $`S=\mathrm{exp}_{\mathrm{e}K}(\mathrm{N}_{\mathrm{e}K}(H_1\mathrm{e}K))`$ is polar and $`[\nu ,\nu ]\text{h}`$, where $`\nu \text{m}_1\text{p}`$ is a normal space to a principal orbit of the slice representation of $`HK`$ on $`\text{m}_1\text{p}`$.
## 12. Subactions of cohomogeneity one and transitive actions
To finish the proof of our classification result it remains to consider subactions of cohomogeneity one and transitive actions. Note that polar actions on the real Grassmannians $`𝔾_2(^n)`$ of rank two were completely classified in and we will not consider any spaces locally isometric to them here. We may also ignore all actions with a fixed point, since they are known to be hyperpolar by Corollary 6.2, see also , Theorem 2.2.
###### Proof of Theorems 1 and 2.
Let $`G`$ be a connected simple compact Lie group and let $`K`$ be a symmetric subgroup such that $`\mathrm{rk}(G/K)2`$. Assume the closed connected subgroup $`HG`$ acts polarly on $`M=G/K`$.
We have already completed the classification in the case where $`G`$ is an exceptional Lie group, hence it remains the case where $`G`$ is one of the classical Lie groups $`\mathrm{SO}(n)`$, $`n7`$, $`\mathrm{SU}(n)`$, $`n3`$, or $`\mathrm{Sp}(n)`$, $`n2`$. Then $`H`$ is contained in one of the maximal connected subgroups of $`G`$ as described in Propositions 3.6, 3.7, and 3.8. Thus at least one of the following holds:
* $`H`$ is contained in a maximal tensor product subgroup (8.2) of $`G`$.
* $`H`$ is contained in maximal connected simple irreducible subgroup of $`G`$.
* $`H`$ is contained in a symmetric subgroup of $`G`$.
If the first possibility holds, then the result follows from Propositions 8.1 and 8.3 except if $`M`$ is a Grassmannian and the tensor product subgroup acts with cohomogeneity one on $`M`$. The second possibility was studied in Section 9, except for subactions of cohomogeneity one or transitive actions. In the case where $`H`$ is contained in symmetric subgroup $`H_1`$ of $`G`$, i.e. the $`H`$-action on $`M`$ is a subaction of a Hermann action; the result follows from Theorem 7.3, under the assumption that the cohomogeneity of the $`H_1`$-action on $`M`$ is $`2`$. Thus it remains the case where $`H`$ is a proper closed connected subgroup of $`H_1G`$ such that $`H_1`$ acts on $`M`$ with cohomogeneity $`1`$.
It will turn out that all polar actions on $`G/K`$ are hyperpolar, hence it follows from Corollary 2.12 of that the sections are embedded submanifolds.
#### Subactions of “exceptional” cohomogeneity one and transitive actions
Let us first consider the case where the $`H_1`$-action on $`M=G/K`$ is not of Hermann type. These cohomogeneity one and transitive actions were determined in , Theorem A and , respectively. We only consider the cases where $`G/K`$ is symmetric of rank $`2`$; these actions are given in Table 9.
We examine these actions case by case. Assume first $`H\mathrm{G}_2`$ is acting on $`𝔾_2(^7)`$ or $`𝔾_3(^7)`$. Then $`H`$ is contained in one of the maximal connected subgroups (10.2). Under the $`7`$-dimensional irreducible orthogonal representation of $`\mathrm{G}_2`$, the first two groups $`\mathrm{SO}(4)`$ and $`\mathrm{SU}(3)`$ are mapped to reducible subgroups of $`\mathrm{SO}(7)`$, thus the $`H`$-action is in this case a subaction of a Hermann action. If $`H`$ is contained in $`\mathrm{SU}(3)`$, then the $`H`$-action is a subaction of a cohomogeneity one Hermann action of type BD I-I and will be treated on page 12. The third group can be excluded by Lemma 3.3.
Let us now consider subgroups of $`\mathrm{Spin}(9)`$, acting on $`𝔾_2(^{16})`$. The maximal connected subgroups of $`\mathrm{SO}(9)`$ are, see Proposition 3.6 and .
$`\mathrm{SO}(8),\mathrm{SO}(7)\times \mathrm{SO}(2),\mathrm{SO}(6)\times \mathrm{SO}(3),\mathrm{SO}(5)\times \mathrm{SO}(4),\mathrm{SO}(3)\mathrm{SO}(3),\mathrm{A}_1^{60}.`$
Let $`H_1\mathrm{Spin}(9)`$ such that $`\pi (H_1)\mathrm{SO}(9)`$ is one of the above, where $`\pi :\mathrm{Spin}(9)\mathrm{SO}(9)`$ is the double cover. We need to consider the image of $`H_1`$ under the spin representation $`\delta :\mathrm{Spin}(9)\mathrm{SO}(16)`$. We have $`\delta (\mathrm{Spin}(8))\mathrm{SO}(8)\times \mathrm{SO}(8)`$ and $`\delta (\mathrm{Spin}(7)\mathrm{Spin}(2))\mathrm{U}(8)`$, thus any subgroups of these are contained in symmetric subgroups of $`\mathrm{SO}(16)`$; the remaining subgroups can be excluded by a dimension count, see Lemma 3.3.
We do not need to consider subactions of the $`\mathrm{Sp}(n)\mathrm{Sp}(1)`$-action on $`𝔾_2(^{4n})=\mathrm{SO}(4n)/\mathrm{SO}(4n2)\times \mathrm{SO}(2)`$, since polar actions on these spaces have been completely classified by Podestà and Thorbergsson .
The last item in Table 9 was treated in Section 10.
#### Subactions of cohomogeneity one and transitive Hermann actions
It now remains to study subactions of cohomogeneity one and transitive Hermann actions. These Hermann actions are listed in Table 10. The column marked with $`\frac{G^{\sigma \tau }}{G^\sigma G^\tau }`$ indicates (the local isometry type) of the symmetric space $`G^{\sigma \tau }/G^\sigma G^\tau `$ whose isotropy representation is equivalent to one slice representation of the $`G^\sigma `$-action on $`G/G^\tau `$ (and of the $`G^\tau `$-action on $`G/G^\sigma `$) by Lemma 11.1. (The presentation may be non-effective, in particular for transitive actions the space $`G^{\sigma \tau }/G^\sigma G^\tau `$ is a noneffective presentation of a zero-dimensional space.) We only have to consider actions on symmetric spaces of rank $`2`$.
#### A III-I
Consider the action of $`H_1=\mathrm{S}(\mathrm{U}(p)\times \mathrm{U}(1))`$ on $`G/K=\mathrm{SU}(p+1)/\mathrm{SO}(p+1)`$, $`p2`$. Assume a closed connected subgroup $`HH_1`$ acts polarly on $`G/K`$. Then $`H`$ is contained in some maximal connected subgroup $`H_2`$ of $`H_1`$. By Theorem 2.1 of , either $`H_2=\mathrm{SU}(p)`$ or $`H_2=\mathrm{S}((H_2^{}\mathrm{U}(1))\times \mathrm{U}(1))`$ where $`H_2^{}`$ is a maximal connected subgroup of $`\mathrm{SU}(p)`$, see Proposition 3.7. In the case of the $`\mathrm{SU}(p)`$-action on $`G/K`$, an explicit calculation shows that the normal space to the orbit at $`\mathrm{e}K`$ is not a Lie triple system, thus the $`\mathrm{SU}(p)`$-action on $`G/K`$ is non-polar by Proposition 4.1. Thus we may restrict our attention to the second case, where we may further assume that $`H_2^{}\mathrm{SU}(p)`$ is irreducible, since otherwise the $`H`$-action is a subaction of a Hermann action whose cohomogeneity is $`2`$. Assume first $`H_2^{}=\mathrm{SO}(p)`$. Then the $`H`$-action on $`G/K`$ is non-polar by Corollary 6.2, since the $`H_2`$-action has a one-dimensional orbit. Now assume $`H_2^{}=\mathrm{Sp}(p/2)`$, $`p4`$, then one isotropy group is $`\mathrm{U}(p/2)`$ and the slice representation is the adjoint representation of $`\mathrm{SU}(p/2)`$ plus the standard representation of $`\mathrm{U}(p/2)`$, see Table 5, which is non-polar and polarity minimal . The normal space $`\mathrm{N}_{\mathrm{e}K}H_2\mathrm{e}K`$ contains a Lie triple system corresponding to an irreducible symmetric space of higher rank, thus the $`H_2`$-action on $`G/K`$ is non-polar and polarity minimal by Lemma 6.4 (iii). Subgroups of $`\mathrm{S}(\mathrm{U}(p/\mathrm{})\mathrm{U}(\mathrm{}))\times \mathrm{U}(1))`$ are excluded by Proposition 8.2, simple irreducible maximal connected subgroups $`H_2\mathrm{SU}(p)`$ by Proposition 9.2 (iii).
#### A III-II
Let $`H`$ be a closed connected subgroup of $`H_1=\mathrm{S}(\mathrm{U}(2nk)\times \mathrm{U}(k))`$, $`k=1,2,3`$ acting on $`G/K=\mathrm{SU}(2n)/\mathrm{Sp}(n)`$, $`n3`$. It is well known that the action of $`H=\mathrm{SU}(2n1)`$ is transitive on $`G/K`$.
We will now study actions of closed connected subgroups $`H`$ in $`H_2=\mathrm{S}((H_2^{}\mathrm{U}(1))\times \mathrm{U}(k))`$. The cases where $`H_2^{}\mathrm{SU}(2nk)`$ is a simple irreducible or tensor product subgroup are excluded by Lemma 9.2 and Proposition 8.2. Thus it remains to consider the case where $`H_2^{}`$ is a symmetric subgroup of $`\mathrm{SU}(2nk)`$. Assume $`H_2^{}=\mathrm{SO}(2nk)`$; if $`k=1`$ then $`\mathrm{SO}(2n1)`$ acts on the symmetric space $`M=\mathrm{SU}(2n)/\mathrm{Sp}(n)`$, homogeneously presented as $`M=\mathrm{SU}(2n1)/\mathrm{Sp}(n1)`$; an isotropy subgroup of this action is $`H_2K=\mathrm{U}(n1)`$, its slice representation is equivalent to the adjoint representation of $`\mathrm{SU}(n1)`$ plus the standard representation on $`^{n1}=^{2n2}`$, see Table 5. This representation is non-polar and polarity minimal by Proposition 6.3 and hence the $`H_2`$-action on $`G/K`$ is non-polar and polarity minimal by Lemma 6.4 (iii), since the normal space contains an irreducible Lie triple system of higher rank. If $`k=2`$ or $`k=3`$, then a slice representation of the $`H_2`$-action on $`G/K`$ contains a module equivalent to the isotropy representation of $`\mathrm{P}^{n1}`$ or $`\mathrm{P}^{n2}`$ restricted to $`\mathrm{U}(n1)\times \mathrm{Sp}(1)`$ or $`\mathrm{U}(n2)\times \mathrm{Sp}(1)`$, respectively. This representation contains two equivalent modules and the $`H`$-action is thus non-polar and polarity minimal by Lemma 6.4 (iii), except if $`n=k=3`$, a case which can be handled by explicit calculations using the criterion in Proposition 4.1.
If $`k=2`$ and $`H_2^{}=\mathrm{Sp}(n1)`$, then the $`H_2`$-action has a one-dimensional orbit and the action of any closed subgroup $`HH_2`$ on $`G/K`$ is non-polar by Corollary 6.2. Thus we are left with the case where $`H_2^{}=\mathrm{U}(2nk\mathrm{})\times \mathrm{U}(\mathrm{})`$. We may assume $`k+\mathrm{}3`$ since otherwise the $`H`$-action on $`G/K`$ is a subaction of a Hermann action of cohomogeneity $`2`$, which were already treated in Section 7. If $`k=\mathrm{}=1`$ then we obtain the cohomogeneity one actions of $`H=\mathrm{S}(\mathrm{U}(2n2)\times \mathrm{U}(1)\times \mathrm{U}(1))`$ and $`\mathrm{S}(\mathrm{U}(2n2)\times \mathrm{U}(1))`$ on $`G/K`$, we have already seen that no further closed proper subgroup of these groups acts polarly. In case $`k+\mathrm{}=3`$ we have to consider closed connected subgroups $`H`$ of $`H_2=\mathrm{S}(\mathrm{U}(2n3)\times (H_2^{\prime \prime }\mathrm{U}(1)))`$, where $`H_2^{\prime \prime }\mathrm{SU}(3)`$ is a maximal connected subgroup. Since there are a number of subgroups $`HH_2`$ acting with cohomogeneity two (in these cases all slice representations are polar), we have to exclude them by explicit calculations using Proposition 4.1.
In case $`n=k=3`$ there are additional maximal connected subgroups, i.e.
$$H_2=\{(zA,z^1\alpha (A))A\mathrm{SU}(3),z,\left|z\right|=1\}\mathrm{S}(\mathrm{U}(3)\times \mathrm{U}(3)),$$
where $`\alpha \mathrm{Aut}(\mathrm{SU}(3))`$, see Theorem 2.1 of . If $`\alpha `$ is an outer automorphism, e.g. given by complex conjugation, then the $`H_2`$-action has a one-dimensional orbit and is non-polar and polarity minimal by Corollary 6.2. If $`\alpha `$ is an inner automorphism, then a stabilizer component is $`\mathrm{U}(1)\mathrm{SO}(3)`$, the $`9`$-dimensional slice representation is equivalent to $`^3^1^3`$ hence non-polar and polarity minimal by Proposition 6.3 (iii). Thus the $`H_2`$-action on $`G/K`$ is non-polar and polarity minimal by Lemma 6.4.
#### A II-III
Consider the action of $`H_1=\mathrm{Sp}(n)`$ on $`G/K=\mathrm{SU}(2n)/\mathrm{S}(\mathrm{U}(2n2)\times \mathrm{U}(2))`$, $`n2`$. Let $`HH_1`$ be a closed connected subgroup acting polarly on $`G/K`$. The $`H_1`$-orbit $`H_1\mathrm{e}K\mathrm{Sp}(n)/\mathrm{Sp}(n1)\times \mathrm{Sp}(1)`$ is totally geodesic and $`H`$ acts polarly on this orbit by Lemma 11.2, the action being non-transitive by . The $`H`$-action on $`H_1\mathrm{e}K`$ has a singular orbit by Corollary 5.3 and we may assume by conjugation of $`H`$ in $`H_1`$ that $`\mathrm{e}K`$ lies in a singular orbit. From Table 10 we read off that the slice representation of the $`H_1`$-action on $`G/K`$ is equivalent to the isotropy representation of the symmetric space $`H_1\mathrm{e}K=H_1/H_1K`$. Thus the nontrivial slice representation of the $`H`$-action on $`H_1\mathrm{e}K`$, which is a submodule of the isotropy representation of $`H_1/H_1K`$, also occurs as a submodule of the slice representation of $`H_1`$ on $`G/K`$ restricted to $`HK`$. We conclude that the slice representation of the $`H`$-action on $`G/K`$ contains two nontrivial equivalent modules and is hence non-polar by , Lemma 2.9.
Consider the action of $`H_1=\mathrm{Sp}(n)`$ on $`G/K=\mathrm{SU}(2n)/\mathrm{S}(\mathrm{U}(2n3)\times \mathrm{U}(3))`$, $`n3`$. Let $`HH_1`$ be a closed connected subgroup acting polarly on $`G/K`$. Then $`H`$ is contained in a maximal connected subgroup $`H_2`$ of $`H_1=\mathrm{Sp}(n)`$. We may assume that $`H_2`$ is irreducible, since otherwise the $`H`$-action is a subaction of a Hermann action with cohomogeneity $`2`$. If $`H_2=\mathrm{U}(n)`$ then $`H_2`$ is contained after conjugation in $`\mathrm{S}(\mathrm{U}(n)\times \mathrm{U}(n))\mathrm{SU}(2n)`$ and the $`H`$-action is also a subaction of a Hermann action of cohomogeneity $`2`$. The actions of maximal connected subgroups of type $`\mathrm{SO}(q)\mathrm{Sp}(p/q)`$ have been treated in Proposition 8.4. The actions of simple irreducible subgroups $`\rho (H)`$, where $`\rho :H\mathrm{Sp}(n)`$ is an irreducible representation of quaternionic type, have been excluded in Section 9.
#### A III-III
Consider the action of $`H_1=\mathrm{S}(\mathrm{U}(a+b)\times \mathrm{U}(1))`$ on the complex Grassmannian $`G/K=\mathrm{SU}(a+b+1)/\mathrm{S}(\mathrm{U}(a)\times \mathrm{U}(b+1))`$, $`ab1`$, $`a+b3`$. Assume $`HH_1`$ is a closed connected subgroup acting polarly on $`G/K`$. First note that the action of $`\mathrm{SU}(a+b)`$ on $`G/K`$ is orbit equivalent to the $`H_1`$-action. Now assume $`HH_2=\mathrm{S}((H_2^{}\mathrm{U}(1))\times \mathrm{U}(1))`$ is a closed connected subgroup acting polarly on $`G/K`$, where $`H_2^{}\mathrm{SU}(a+b)`$ is a maximal connected subgroup. We may assume that the standard representation of $`\mathrm{SU}(a+b)`$ restricted to $`H_2^{}`$ acts irreducibly on $`^{a+b}`$, since otherwise $`H_2`$ is contained in a subgroup of $`G`$ conjugate to $`\mathrm{S}(\mathrm{U}(k)\times \mathrm{U}(a+b+1k))`$ for $`2ka+b1`$ and the $`H_2`$-action on $`G/K`$ is a subaction of a Hermann action of cohomogeneity $`2`$, which have already been treated in Section 7. It follows from Lemma 11.2 that the orbit $`H_1\mathrm{e}K`$ is a totally geodesic submanifold of $`G/K`$ isometric to $`\mathrm{SU}(a+b)/\mathrm{S}(\mathrm{U}(a)\times \mathrm{U}(b))`$ on which $`H`$ acts polarly.
Assume first that $`H_2^{}`$ is an irreducible symmetric subgroup of $`\mathrm{SU}(a+b)`$, hence conjugate to either $`\mathrm{SO}(a+b)`$ or $`\mathrm{Sp}(\frac{a+b}{2})`$. However, in the first case the isotropy subgroup of the $`H_2`$-action at $`\mathrm{e}K`$ is $`\mathrm{S}(\mathrm{O}(a)\times \mathrm{O}(b))\mathrm{U}(1)`$ and its slice representation is equivalent to $`(^a^b)(^a^1)`$, thus it is non-polar and polarity minimal by Lemma 6.3. It follows from Table 10 that the normal space contains a Lie triple system corresponding to a totally geodesic submanifold isometric to $`\mathrm{P}^a`$, thus the $`H_2`$-action is non-polar and polarity minimal by Lemma 6.4 (iii), since $`a2`$. Let us now consider the case where $`a+b`$ is even and $`H_2^{}`$ is conjugate to $`\mathrm{Sp}(\frac{a+b}{2})`$, hence $`a+b4`$. The group $`H`$ acts polarly on the totally geodesic $`H_1`$-orbit $`H_1\mathrm{e}K\mathrm{SU}(a+b)/\mathrm{S}(\mathrm{U}(a)\times \mathrm{U}(b))`$, which is of rank $`b`$. If $`b2`$, then the reducible slice representation is non-polar and polarity minimal by , and the $`H`$-action is non-polar and polarity minimal by Lemma 6.4. If $`b=1`$, then $`H_2`$ acts transitively on the orbit $`H_1\mathrm{e}K`$, but an explicit calculation using Proposition 4.1 shows that the $`H_2`$-action on $`G/K`$ is non-polar. Let $`HH_2`$ be a proper closed subgroup, then $`H`$ acts non-transitively on $`H_1\mathrm{e}K`$ by . If the $`H`$-action on $`G/K`$ is polar, then also the $`H`$-action restricted to $`H_1\mathrm{e}K\mathrm{P}^a`$ is polar by Lemma 4.2 and it has a singular orbit $`Hp`$ by Corollary 5.3. The normal space $`\text{p}\text{m}_1`$, see Lemma 11.1, of the $`H_1`$-action on $`G/K`$ contains a submodule which is equivalent to the slice representation at $`pH_1\mathrm{e}K`$ of the $`H`$-action on $`H_1\mathrm{e}K`$ after a $`\mathrm{U}(1)`$-factor is removed from both representations. Since both modules belong to the polar slice representation of the $`H`$-action on $`G/K`$, it follows from that $`H`$ is at most three-dimensional, a contradiction with Proposition 5.5.
Now assume $`H_2^{}`$ is a non-symmetric irreducible maximal connected subgroup of $`\mathrm{Sp}(a+b)`$. It follows from what we have shown so far that this can only happen if $`\mathrm{rk}(H_1\mathrm{e}K)=b=1`$. Assume $`H_2^{}=\mathrm{SO}(p)\mathrm{Sp}(q)`$, then the $`H_2`$-action on $`H_1\mathrm{e}K`$ is non-polar and polarity minimal by Proposition 8.3. If $`H_2^{}`$ is a simple irreducible maximal connected subgroup of $`\mathrm{Sp}(a+b)`$, then it follows from the results of Section 9 that the action of $`H_2`$ on $`G/K`$ is non-polar and polarity minimal, since if the action of $`H_2^{}`$ on $`𝔾_k(^{a+b})`$ for $`2ka+b2`$ is excluded by Proposition 5.6, then also the action of $`H_2`$ on $`𝔾_k(^{a+b+1})`$ is excluded by a dimension count.
#### BD I-I
Let $`H_1=\mathrm{SO}(a+b)`$, $`G/K=\mathrm{SO}(a+b+1)/\mathrm{SO}(a)\times \mathrm{SO}(b+1)`$, $`a+b6`$, $`ab1`$. Assume $`HH_1`$ is a closed connected subgroup acting polarly on $`G/K`$. Without loss of generality we may assume that $`H\mathrm{SO}(a+b)`$ acts irreducibly on $`^{a+b}`$, since otherwise the $`H`$-action on $`G/K`$ is a subaction of a Hermann action of cohomogeneity $`2`$, see Section 7. By Lemma 11.2, the $`H_1`$-orbit $`H_1\mathrm{e}K`$ is a totally geodesic submanifold isometric to $`\mathrm{SO}(a+b)/\mathrm{S}(\mathrm{O}(a)\times \mathrm{O}(b))`$ and $`H`$ acts polarly on $`H_1\mathrm{e}K`$ (the action may be transitive).
We do not need to consider the case where $`\mathrm{rk}(H_1\mathrm{e}K)=b=1`$, since polar actions on $`𝔾_2(^{a+2})`$ were completely classified in , hence we may assume $`b2`$. Let us first consider the cases where $`HH_1`$ is not a symmetric subgroup. Then it follows from what we have shown so far and that we have one of the following:
* $`H=\mathrm{G}_2,a+b=7,b=2,\mathrm{\hspace{0.17em}3}`$;
* $`H=\mathrm{Spin}(7),a+b=8,b=2,\mathrm{\hspace{0.17em}3},\mathrm{\hspace{0.17em}4}`$;
* $`H=\mathrm{Spin}(9),a+b=16,b=2`$;
* $`H=\mathrm{Sp}(n)\mathrm{Sp}(1),a+b=4n,b=2`$;
* $`H=\mathrm{U}(n),a+b=2n`$.
In the case of the $`\mathrm{G}_2`$-actions, an explicit calculation using Proposition 4.1 shows that the actions are non-polar; subgroups of $`\mathrm{G}_2`$, see (10.2), are either reducible or are excluded by a dimension count. The actions of $`\mathrm{Spin}(7)`$ are orbit equivalent to the $`\mathrm{SO}(8)`$-action in case $`b=2,\mathrm{\hspace{0.17em}3}`$, in case $`b=4`$ the action can be shown to be non-polar by an explicit calculation. Subgroups of $`\mathrm{Spin}(7)`$ are either contained in groups treated below or ruled out by a dimension count. The actions of $`\mathrm{Spin}(9)`$ and $`\mathrm{Sp}(n)\mathrm{Sp}(1)`$ can be excluded by replacing $`K`$ with the conjugate subgroup $`K^{}=\mathrm{SO}(3)\times \mathrm{SO}(a)`$, the actions on $`H_1\mathrm{e}K^{}`$ have already been shown to be non-polar and polarity minimal. Assume now $`H=\mathrm{U}(\frac{a+b}{2})`$ . The slice representation of the $`H`$-action at $`\mathrm{e}K`$, as can be seen from Tables 5 and 10, contains a module equivalent to the representation of $`\mathrm{U}\left(\frac{a}{2}\right)\times \mathrm{U}\left(\frac{b}{2}\right)`$ on $`^{\frac{a}{2}}^{\frac{b}{2}}^{\frac{a}{2}}`$, which is non-polar and polarity minimal by Lemma 6.3 since $`\frac{a}{2}2`$.
#### C I-II
Let $`H`$ be closed connected subgroup of $`H_1=\mathrm{Sp}(p)\times \mathrm{Sp}(1)`$ acting polarly on $`G/K=\mathrm{Sp}(p+1)/\mathrm{U}(p+1)`$, $`p2`$. We first observe that the actions of $`\mathrm{Sp}(p)`$ and $`\mathrm{Sp}(p)\times \mathrm{U}(1)`$ are not orbit equivalent to the $`H_1`$-action; since the normal space at a regular orbit is not a Lie triple system, these actions are non-polar by Proposition 4.1. Now assume $`HH_2=H_2^{}\times \mathrm{Sp}(1)`$, where $`H_2^{}`$ is a maximal connected subgroup of $`\mathrm{Sp}(n)`$. We may assume $`H_2^{}\mathrm{Sp}(n)`$ acts irreducibly on $`^n`$, since otherwise $`H_2`$ is a subgroup of $`\mathrm{Sp}(p+1k)\times \mathrm{Sp}(k)`$, $`2kp1`$, see Section 7. Consider the action of $`H_2=\mathrm{U}(p)\times \mathrm{Sp}(1)`$ on $`G/K`$, then the slice representation of the isotropy subgroup $`\mathrm{U}(p)\times \mathrm{U}(1)`$ is equivalent to the isotropy representation of $`\mathrm{Sp}(p)/\mathrm{U}(p)`$ plus $`^p^1`$. This representation is non-polar and the action of $`H_2`$ on $`G/K`$ is polarity-minimal by Lemma 6.4 (iii). If $`H_2=\rho (H_2^{})\times \mathrm{Sp}(1)`$, where $`\rho :H_2^{}\mathrm{Sp}(p)`$ is an irreducible representation of the simple compact Lie group $`H`$, then the $`H`$-action on $`G/K`$ is non-polar by Lemma 9.2. Tensor product subgroups $`H_2^{}`$ have been excluded in Proposition 8.2.
#### C II-II
Let $`H_1=\mathrm{Sp}(a+b)\times \mathrm{Sp}(1)`$ and let $`G/K=\mathrm{Sp}(a+b+1)/\mathrm{Sp}(a)\times \mathrm{Sp}(b+1)`$, $`ab1`$, $`a+b3`$. First observe that the action of $`\mathrm{Sp}(a+b)`$ on $`G/K`$ is orbit equivalent to the $`H_1`$-action. Now assume $`HH_2=H_2^{}\times \mathrm{Sp}(1)`$ is a closed connected subgroup acting polarly on $`G/K`$, where $`H_2^{}\mathrm{Sp}(a+b)`$ is a maximal connected subgroup. We may assume that $`H_2^{}\mathrm{Sp}(a+b)`$ acts irreducibly on $`^{a+b}`$, since otherwise the $`H`$-action on $`G/K`$ is a subaction of a Hermann action of cohomogeneity $`2`$, which were examined in Section 7. It follows from Lemma 11.2 that the $`H_1`$-orbit $`H_1\mathrm{e}K`$ is a totally geodesic submanifold isometric to $`\mathrm{Sp}(a+b)/\mathrm{Sp}(a)\times \mathrm{Sp}(b)`$ on which $`H`$ acts polarly.
Assume $`H_2^{}`$ is an irreducible symmetric subgroup of $`\mathrm{Sp}(a+b)`$, hence conjugate to $`\mathrm{U}(a+b)`$. However, the isotropy subgroup of the $`H_2`$-action at $`\mathrm{e}K`$ is $`\mathrm{U}(a)\times \mathrm{U}(b)\times \mathrm{Sp}(1)`$, its slice representation contains two equivalent modules, thus it is non-polar and polarity minimal by Lemma 6.3, parts (ii) and (iii). The normal space contains a Lie triple system corresponding to a totally geodesic submanifold isometric to $`\mathrm{P}^a`$, hence the $`H_2`$-action is non-polar and polarity minimal by Lemma 6.4 (iii).
Now assume $`H_2^{}`$ is a non-symmetric irreducible maximal connected subgroup of $`\mathrm{Sp}(a+b)`$ acting polarly on $`H_1\mathrm{e}K=\mathrm{Sp}(a+b)/\mathrm{Sp}(a)\times \mathrm{Sp}(b)`$. It follows from what we have shown so far that this can only happen if $`\mathrm{rk}(H_1\mathrm{e}K)=1`$, i.e. $`b=1`$. Assume $`H_2^{}=\mathrm{SO}(p)\mathrm{Sp}(q)`$, then the $`H_2`$-action on $`H_1\mathrm{e}K`$ is non-polar and polarity minimal by Proposition 8.3. If $`H_2^{}`$ is a simple irreducible subgroup of $`\mathrm{Sp}(a+b)`$, then it follows from the results of Section 9 that the action of $`H_2`$ on $`G/K`$ is non-polar and polarity minimal, since if the action of $`H_2^{}`$ on $`𝔾_k(^{a+b})`$ for $`2ka+b2`$ is excluded by Proposition 5.6, then so is the action of $`H_2`$ on $`𝔾_k(^{a+b+1})`$.
#### D I-III
Let $`H`$ be a closed connected subgroup of $`\mathrm{SO}(nk)\times \mathrm{SO}(k)`$, $`k=1,2,3`$ acting on $`G/K=\mathrm{SO}(2n)/\mathrm{U}(n)`$, $`n3`$. We first study actions of closed connected subgroups $`H`$ in $`H_2=H_2^{}\times \mathrm{SO}(k)`$. The cases where $`H_2^{}\mathrm{SO}(2nk)`$ is a simple irreducible or tensor product subgroup were excluded by Lemma 9.2 and Lemma 8.2. Let us consider the case where $`H_2^{}`$ is a symmetric subgroup of $`\mathrm{SO}(2nk)`$. If $`k=2`$ and $`H_2^{}=\mathrm{U}(n1)`$, then the $`H`$-action has a fixed point. Thus it remains the case where $`H_2^{}=\mathrm{SO}(2nk\mathrm{})\times \mathrm{SO}(\mathrm{})`$. We may assume $`k+\mathrm{}3`$ since otherwise the $`H`$-action on $`G/K`$ is a subaction of a Hermann action of cohomogeneity $`2`$, see Section 7. If $`k=\mathrm{}=1`$ then we obtain the cohomogeneity one action of $`H=\mathrm{SO}(2n2)`$ on $`G/K`$, we have already seen that no closed proper subgroup of this group acts polarly. In the case where $`k+\mathrm{}=3`$, an explicit calculation using Proposition 4.1 shows that the actions of $`\mathrm{SO}(2n3)`$ and $`\mathrm{SO}(2n3)\times \mathrm{SO}(2)`$ on $`G/K`$ are non-polar; we have already excluded any closed subgroups of these two groups.
In case $`n=3`$ there is an additional maximal connected subgroup, i.e. $`\mathrm{\Delta }\mathrm{SO}(3)=\{(g,g)g\mathrm{SO}(3)\}\mathrm{SO}(3)\times \mathrm{SO}(3)`$, but this action has a fixed point.
#### D III-I
Let $`H`$ be a closed connected subgroup of $`H_1=\mathrm{U}(n)`$ acting polarly on the symmetric space $`G/K=\mathrm{SO}(2n)/\mathrm{SO}(n3)\times \mathrm{SO}(3)`$, $`n3`$. It follows from Theorem 2.1 in that the conjugacy classes of maximal connected subgroups $`H_2`$ in $`\mathrm{U}(n)`$ are exhausted by $`\mathrm{SU}(n)`$, and $`H_2^{}\mathrm{U}(1)`$ where $`H_2^{}`$ runs through the maximal connected subgroups of $`\mathrm{SU}(n)`$; see Proposition 3.7. We observe first that the action of $`\mathrm{SU}(n)`$ is orbit equivalent to the $`\mathrm{U}(n)`$-action. Now assume $`HH_2^{}\mathrm{U}(1)`$. We do not need to consider reducible subgroups $`H_2^{}`$ since they lead to subactions of Hermann actions with cohomogeneity $`2`$, which were treated in Section 7. Also, if $`H_2^{}=\mathrm{SO}(n)`$, then the same argument as in Proposition 8.3 shows that the action is non-polar and polarity minimal. If $`H_2^{}=\mathrm{Sp}(n/2)`$, then $`H_2`$ is contained in $`\mathrm{Sp}(n/2)\mathrm{Sp}(1)`$ and the action of any closed subgroup $`HH_2`$ was shown not to be polar in Proposition 8.3. Assume now $`H_2^{}`$ is a tensor product subgroup $`\mathrm{SU}(p)\mathrm{SU}(n/p)`$, then the action of $`H`$ is non-polar according to Proposition 8.4. Finally, let $`H_2^{}`$ be given by an irreducible representation $`\rho :H_2^{}\mathrm{SU}(n)`$ where $`H_2^{}`$ is a simple compact Lie group; these actions were excluded in Section 9.
#### E II-IV
Let $`HH_1=\mathrm{SU}(6)\mathrm{Sp}(1)`$ be a closed connected subgroup acting on $`G/K=\mathrm{E}_6/\mathrm{F}_4`$. Then $`H`$ is contained in one of the maximal connected subgroups of $`\mathrm{SU}(6)\mathrm{Sp}(1)`$. By , we may assume the involutions of $`\mathrm{E}_6`$ corresponding to $`H_1`$ and $`K`$ commute such that the totally geodesic $`H_1`$-orbit through $`\mathrm{e}K`$ is isometric to $`\mathrm{SU}(6)\mathrm{Sp}(1)/\mathrm{Sp}(3)\mathrm{Sp}(1)\mathrm{SU}(6)/\mathrm{Sp}(3)`$, see Lemma 11.2; the slice representation is equivalent to the isotropy representation of $`\mathrm{Sp}(4)/\mathrm{Sp}(3)\mathrm{Sp}(1)`$, this shows that the $`\mathrm{Sp}(1)`$-factor is inessential for the $`H_1`$-action. Now assume $`H_2^{}\mathrm{SU}(6)`$ is a symmetric subgroup and $`HH_2=H_2^{}\mathrm{Sp}(1)`$ is a closed connected subgroup. If $`H_2^{}=\mathrm{Sp}(3)`$ then the $`H`$-action on $`G/K`$ has a fixed point. If $`H_2^{}=\mathrm{SO}(6)`$, then the connected component of an isotropy subgroup is $`\mathrm{U}(3)\times \mathrm{Sp}(1)`$ and its slice representation is equivalent to the adjoint representation of $`\mathrm{SU}(3)`$ plus $`\chi (\mathrm{Sp}(4),\mathrm{Sp}(3)\times \mathrm{Sp}(1))|_{\mathrm{U}(3)\times \mathrm{Sp}(1)}`$, see Tables 5 and 10, it is non-polar and polarity minimal by Proposition 6.3. Hence the $`H_2`$-action on $`G/K`$ is non-polar and polarity minimal by Lemma 6.4 (iii).
Now assume $`H_2^{}`$ is one of the groups $`\mathrm{S}(\mathrm{U}(k)\times \mathrm{U}(6k))`$, $`k=1,2,3`$. If $`k=3`$, then $`(H_2K)_0\mathrm{Sp}(1)\mathrm{U}(1)\mathrm{Sp}(1)\mathrm{Sp}(1)`$ and the slice representation is equivalent to
$$\chi (\mathrm{Sp}(2),\mathrm{Sp}(1)\times \mathrm{Sp}(1))\chi (\mathrm{Sp}(4),\mathrm{Sp}(3)\times \mathrm{Sp}(1))|_{\mathrm{Sp}(1)\mathrm{U}(1)\mathrm{Sp}(1)\mathrm{Sp}(1)},$$
where both modules are polar, but their sum is non-polar, see ; using the results of , it can be directly verified that this representation is polarity minimal and it follows that the $`H_2`$-action on $`G/K`$ is non-polar and polarity minimal by Lemma 6.4. If $`k=2`$, then $`(H_2K)_0\mathrm{Sp}(2)\mathrm{Sp}(1)\mathrm{Sp}(1)`$ and the corresponding slice representation is equivalent to
$$\chi (\mathrm{Sp}(3),\mathrm{Sp}(2)\times \mathrm{Sp}(1))\chi (\mathrm{Sp}(4),\mathrm{Sp}(3)\times \mathrm{Sp}(1))|_{\mathrm{Sp}(2)\mathrm{Sp}(1)\mathrm{Sp}(1)}$$
hence non-polar by and polarity minimal by Proposition 6.3 (ii) and the $`H_2`$-action on $`G/K`$ is non-polar and polarity minimal by Lemma 6.4. Finally, if $`k=1`$ then $`H_2`$ acts transitively on the $`H_1`$-orbit through $`\mathrm{e}K`$, and the group $`(H_2K)_0\mathrm{U}(1)\times \mathrm{Sp}(2)\times \mathrm{Sp}(1)`$ acts polarly on $`\mathrm{exp}_{\mathrm{e}K}(\mathrm{N}_{\mathrm{e}K}(H_2\mathrm{e}K)\mathrm{P}^3`$ by , cf. Lemma 11.1; however, an explicit calculation shows that $`[\nu ,\nu ]⟂̸\text{h}_2`$, where $`\nu \text{m}_1\text{p}`$ is the tangent space to a section of the $`H_2K`$-action and hence the $`H_2`$-action and the orbit equivalent action of $`\mathrm{SU}(5)\mathrm{Sp}(1)`$ are non-polar by Lemma 11.3. The actions of $`\mathrm{S}(\mathrm{U}(1)\times \mathrm{U}(5))`$ and $`\mathrm{S}(\mathrm{U}(1)\times \mathrm{U}(5))\mathrm{U}(1)`$ are non-polar since the slice representations at $`\mathrm{e}K`$ are non-polar. Any other closed subgroups of $`H_2=\mathrm{S}(\mathrm{U}(1)\times \mathrm{U}(5))\mathrm{Sp}(1)`$ are contained in the groups treated above or excluded by Proposition 5.6.
#### E IV-II
Consider now closed connected subgroups $`H`$ of $`H_1=\mathrm{F}_4`$ acting polarly on $`\mathrm{E}_6/\mathrm{SU}(6)\mathrm{Sp}(1)`$. It follows from Lemma 3.3 that $`dim(H)28`$. By , the only closed connected subgroup of $`\mathrm{F}_4`$ of sufficient dimension is $`\mathrm{Spin}(9)`$. By conjugation, the subgroup $`H=\mathrm{Spin}(9)\mathrm{F}_4`$ can be chosen such that the connected component of the isotropy group $`(HK)_0`$ is $`\mathrm{Sp}(2)\mathrm{Sp}(1)^2\mathrm{Spin}(5)\mathrm{Spin}(4)`$, see Table 10. From Table 10 one sees further that the slice representation restricted to $`(HK)_0`$ is equivalent to the isotropy representation of $`\mathrm{Sp}(3)/\mathrm{Sp}(2)\mathrm{Sp}(1)`$ plus the isotropy representation of $`\mathrm{Sp}(4)/\mathrm{Sp}(3)\mathrm{Sp}(1)`$ restricted to $`\mathrm{Sp}(2)\mathrm{Sp}(1)`$, hence non-polar by .
#### E III-IV
Let $`H`$ be a closed connected subgroup of $`H_1=\mathrm{Spin}(10)\mathrm{U}(1)`$ acting polarly on $`G/K=\mathrm{E}_6/\mathrm{F}_4`$. Let us first show that the action of $`\mathrm{Spin}(10)`$ on $`G/K`$ is non-polar. Assume the converse, i.e. the action of $`H=\mathrm{Spin}(10)`$ on $`G/K=\mathrm{E}_6/\mathrm{F}_4`$ is polar. Since this action is of cohomogeneity two, it follows from Proposition 2.8 that the sections are locally isometric to a two-sphere. The $`H_1`$-orbit through $`\mathrm{e}K`$ is totally geodesic by Lemma 11.2 and locally isometric to $`\mathrm{S}^9\times \mathrm{S}^1`$, where the $`\mathrm{S}^1`$-factor is the orbit of the $`\mathrm{U}(1)`$-factor in $`H_1=\mathrm{Spin}(10)\mathrm{U}(1)`$, hence totally geodesic in $`G/K`$. It follows from the decomposition (11.2) that $`\mathrm{T}_{\mathrm{e}K}\mathrm{S}^1\mathrm{T}_{\mathrm{e}K}(H\mathrm{e}K)`$, hence $`\mathrm{T}_{\mathrm{e}K}\mathrm{S}^1\mathrm{T}_{\mathrm{e}K}\mathrm{\Sigma }`$, where $`\mathrm{\Sigma }`$ is a section of the $`H`$-action on $`G/K`$ containing $`\mathrm{e}K`$. Since the Lie algebra of the $`\mathrm{U}(1)`$-factor in $`H_1=\mathrm{Spin}(10)\mathrm{U}(1)`$ is contained in $`\mathrm{T}_{\mathrm{e}K}\mathrm{\Sigma }\text{p}`$, it follows by Proposition 3.2 that this $`\mathrm{U}(1)`$-factor acts on $`\mathrm{\Sigma }`$ as a group of transvections. Now, since this $`\mathrm{U}(1)`$-action commutes with the $`H`$-action, it follows that any two points of $`\mathrm{\Sigma }`$ which lie in the same $`\mathrm{U}(1)`$-orbit are of the same orbit type with respect to the $`H`$-action on $`G/K`$. In particular, all singular orbits of the $`H`$-action on $`G/K`$ intersect $`\mathrm{\Sigma }`$ in the $`\mathrm{U}(1)`$-orbit which is covered by a great circle of $`\stackrel{~}{\mathrm{\Sigma }}\mathrm{S}^2`$, since all reflection hypersurfaces $`\{P_j\}_{jJ}`$ have to be invariant under the $`\mathrm{U}(1)`$-action induced on $`\stackrel{~}{\mathrm{\Sigma }}`$, see Lemma 5.1. However, the $`\mathrm{U}(1)`$-action on $`\mathrm{\Sigma }`$, which is isometric to $`\mathrm{P}^2`$ or $`\mathrm{S}^2`$, has at least one fixed point $`p\mathrm{\Sigma }`$. It follows that this point $`pG/K`$ lies in a regular orbit of the $`H`$-action on $`G/K`$, but is left fixed by the $`\mathrm{U}(1)`$-factor of $`H_1=\mathrm{Spin}(10)\mathrm{U}(1)`$. Hence the connected component of the isotropy subgroup $`(H_1)_p`$ is a subgroup $`L\mathrm{U}(1)`$, where $`L\mathrm{Spin}(10)`$ is $`20`$-dimensional, a contradiction, since the principal isotropy subgroup is $`\mathrm{Spin}(7)`$.
Now we may assume that the group $`H`$ acting polarly on $`G/K`$ is contained in $`H_2^{}\mathrm{U}(1)`$, where $`H_2^{}`$ is a maximal connected subgroup of $`\mathrm{Spin}(10)`$. It follows from Lemma 3.3 that $`dim(H)16`$. This implies that $`H_2^{}`$ is one of the following, see Proposition 3.6:
(12.1)
$$\mathrm{Spin}(9),\mathrm{Spin}(8)\mathrm{SO}(2),\mathrm{Spin}(7)\mathrm{Spin}(3),\mathrm{Spin}(6)\mathrm{Spin}(4),\mathrm{U}(5).$$
The actions of these groups are non-polar and polarity minimal by Lemma 6.4 (iv). If $`H_2^{}=\mathrm{Spin}(9)`$, then the $`H`$-action can also be shown to be non-polar by Corollary 6.2.
#### E IV-III
Assume $`H`$ is a closed connected subgroup of $`\mathrm{F}_4`$ acting polarly on $`G/K=\mathrm{E}_6/\mathrm{Spin}(10)\mathrm{U}(1)`$. Proposition 5.5 implies $`dim(H)28`$. By , the only maximal connected subgroups $`H\mathrm{F}_4`$ of dimension $`28`$ is $`\mathrm{Spin}(9)`$. It follows from Table 10 that the action of $`\mathrm{Spin}(9)`$ leaves a point fixed.
#### F II-I
Let $`H`$ be a closed connected subgroup of $`H_1=\mathrm{Spin}(9)`$ acting polarly on $`G/K=\mathrm{F}_4/\mathrm{Sp}(3)\mathrm{Sp}(1)`$. Lemma 3.3 implies $`dim(H)20`$. The only closed connected subgroups of $`\mathrm{Spin}(9)`$ of dimension $`20`$ are $`\mathrm{Spin}(8)`$, $`\mathrm{Spin}(7)\mathrm{SO}(2)`$ and $`\mathrm{Spin}(7)`$.
The subgroup $`H=\mathrm{Spin}(8)\mathrm{Spin}(9)`$ may be chosen such that $`(H_1K)_0=\mathrm{Sp}(2)\mathrm{Sp}(1)\mathrm{Spin}(5)\mathrm{Spin}(3)`$. The group $`H`$ acts with cohomogeneity one on the orbit $`H_1\mathrm{e}K`$, which is covered by $`\mathrm{Spin}(9)/\mathrm{Spin}(5)\mathrm{Spin}(4)`$. With our choice of the subgroup $`\mathrm{Spin}(8)\mathrm{Spin}(9)`$, the slice representation of the $`H`$-action on $`G/K`$ at $`\mathrm{e}K`$ is equivalent (on the Lie algebra level) to the representation of $`\mathrm{Sp}(2)\mathrm{Sp}(1)\mathrm{Spin}(5)\mathrm{Spin}(3)`$ on $`^2_{}^1^5`$, hence non-polar by .
The action of $`\mathrm{Spin}(7)\mathrm{SO}(2)`$ has an isotropy subgroup whose connected component is isomorphic to $`\mathrm{Spin}(5)\mathrm{SO}(2)\mathrm{SO}(2)\mathrm{Sp}(2)\mathrm{U}(1)\mathrm{U}(1)`$ and whose slice representation is $`^5^2\chi (\mathrm{Sp}(3),\mathrm{Sp}(2)\mathrm{Sp}(1))|_{\mathrm{Sp}(2)\mathrm{U}(1)}`$, hence non-polar . This also shows that the $`\mathrm{Spin}(7)`$-action is non-polar. ∎ |
warning/0506/cond-mat0506807.html | ar5iv | text | # Scaling in critical random Boolean networks
## 1 Introduction
Random Boolean networks are often used as generic models for the dynamics of complex systems of interacting entities, such as social and economic networks, neural networks, and gene or protein interaction networks Kauffman et al. (2003). The simplest and most widely studied of these models was introduced in 1969 by Kauffman Kauffman (1969) as a model for gene regulation. The system consists of $`N`$ nodes, each of which receives input from $`K`$ randomly chosen other nodes. The network is updated synchronously, the state of a node at time step $`t`$ being a Boolean function of the states of the $`K`$ input nodes at the previous time step, $`t1`$. The Boolean updating functions are randomly assigned to every node in the network, and together with the connectivity pattern they define the realization of the network. For any initial condition, the network eventually settles on a periodic attractor. Of special interest are *critical* networks, which lie at the boundary between a frozen phase and a chaotic phase Derrida and Pomeau (1986); Derrida and Stauffer (1986). In the frozen phase, a perturbation at one node propagates during one time step on an average to less than one node, and the attractor lengths remain finite in the limit $`N\mathrm{}`$. In the chaotic phase, the difference between two almost identical states increases exponentially fast, because a perturbation propagates on an average to more than one node during one time step Aldana-Gonzalez et al. (2003).
The nodes of a critical network can be classified according to their dynamics on an attractor. First, there are nodes that are frozen on the same value on every attractor. Such nodes give a constant input to other nodes and are otherwise irrelevant. They form the *frozen core* of the network. Second, there are nodes whose outputs go only to irrelevant nodes. Though they may fluctuate, they are also classified as irrelevant since they act only as slaves to the nodes determining the attractor period. Third, the *relevant nodes* are the nodes whose state is not constant and that control at least one relevant node. These nodes determine completely the number and period of attractors. If only these nodes and the links between them are considered, these nodes form loops with possibly additional links and chains of relevant nodes within and between loops. The recognition of the relevant elements as the only elements influencing the asymptotic dynamics was an important step in understanding the attractors of Kauffman networks. The behavior of the frozen core was first studied by Flyvbjerg Flyvbjerg (1988). Then, in an analytical study of $`K=1`$ networks Flyvbjerg and Kjaer Flyvbjerg and Kjær (1988) introduced the concept of relevant elements (though without using this name). The definition of relevant elements that we are using here was given by Bastolla and Parisi Bastolla and Parisi (1998a, b). They gained insight into the properties of the attractors of the critical networks by using numerical experiments based on the modular structure of the relevant elements. Finally, Socolar and Kauffman Socolar and Kauffman (2003) found numerically that for critical $`K=2`$ networks the mean number of nonfrozen nodes scales as $`N^{2/3}`$, and the mean number of relevant nodes scales as $`N^{1/3}`$. The same result is hidden in the analytical work on attractor numbers by Samuelsson and Troein Samuelsson and Troein (2003), as was shown in Drossel (2005).
In this work, we go a step further by deriving these power laws analytically for a more general class of networks, and by showing the asymptotic probability distribution of nonfrozen and relevant nodes in terms of scaling variables. We also obtain results for the number of nonfrozen nodes with two nonfrozen inputs, and for the number of relevant nodes with two relevant inputs. The outline of this paper is the following. In the next section we define the class of networks that we are investigating. In Section 3, we introduce a stochastic process that determines the frozen core of the network starting from the nodes whose outputs are entirely independent of their inputs. Then, in Section 4, we analyze the Langevin and Fokker-Planck equations that correspond to this stochastic process and that lead to the scaling behavior of the number of nonfrozen nodes. In order to identify the relevant nodes among the nonfrozen ones, we introduce in Section 5 another stochastic process. This process also enables us to find their scaling behavior. Finally, we discuss in the last section the implications of our results.
## 2 Critical $`K=2`$ networks
The networks we are studying in this paper are the $`K=2`$ critical networks. In these networks each node has 2 randomly chosen inputs. The 16 possible update functions are shown in table 1.
The update functions fall into four classes Aldana-Gonzalez et al. (2003). In the first class, denoted by $``$, are the frozen functions, where the output is fixed irrespectively of the input. The class $`𝒞_1`$ contains those functions that depend only on one of the two inputs, but not on the other one. The class $`𝒞_2`$ contains the remaining canalizing functions, where one state of each input fixes the output. The class $``$ contains the two reversible update functions, where the output is changed whenever one of the inputs is changed. Critical networks are those where a change in one node propagates to one other node on an average. A change propagates with probability $`1/2`$ to a node that has a canalizing update function $`𝒞_1`$ or $`𝒞_2`$, with probability zero to a node that has a frozen update function, and with probability 1 to a node that has a reversible update function. Consequently, if the frozen and reversible update functions are chosen with equal probability, the network is critical. Usually, only those models are considered where all 16 update functions receive equal weight. We here consider the larger set of models where the frozen and reversible update functions are chosen with equal (and nonzero) probability, and where the remaining probability is divided between the $`𝒞_1`$ and $`𝒞_2`$ functions. Those networks that contain only $`𝒞_1`$ functions are different from the remaining ones. Since all nodes respond only to one input, the link to the second input can be cut, and we are left with a critical $`K=1`$ network, which was already discussed in Flyvbjerg and Kjær (1988); Drossel et al. (2005); Drossel (2005) and will not be discussed here. All the other models, where the weight of the $`𝒞_1`$ functions is smaller than 1, fall into the same class Drossel (2005). The treatment presented in the following, is based on the existence of nodes with frozen functions, and it therefore applies to all critical models with a nonzero fraction of frozen functions. Networks with only canalyzing functions have to be discussed separately.
Let $`N_f`$ be the number of nodes with a frozen function, $`N_r`$ the number of nodes with a reversible function and $`N_{c_1}`$ and $`N_{c_2}`$ the number of nodes with a $`𝒞_1`$ and a $`𝒞_2`$ function. We define the systems we are going to consider through parameters $`\alpha =N_{c_1}/N`$, $`\beta =N_r/N=N_f/N`$, $`\gamma =N_{c_2}/N`$. These parameters give the fraction of each type of nodes in the network. In the next two sections, we determine the properties of the frozen core in the large $`N`$ limit by starting from the nodes with a frozen function.
## 3 A stochastic process that leads to the frozen core
We consider the ensemble of all networks of size $`N`$ and with fixed parameters $`\alpha ,\beta ,\gamma `$. All nodes with a frozen update function are certainly part of the frozen core. We now construct the frozen core by determining stepwise all those nodes that become frozen due to the influence of a frozen node. In the language of Socolar and Kauffman (2003), this process determines the “clamped” nodes. Initially, we place the nodes in four containers labelled $``$, $`𝒞_1`$, $`𝒞_2`$, and $``$. These containers contain $`N_f`$, $`N_{c_1}`$, $`N_{c_2}`$, and $`N_r`$ nodes initially. Since these numbers change during our stochastic process, we denote the initial values as $`N_f^{ini}`$, $`N_{c_1}^{ini}`$, $`N_{c_2}^{ini}`$, and $`N_r^{ini}`$, and the total number of nodes as $`N^{ini}`$. We treat the nodes in container $`𝒞_1`$ as nodes with only one input and with the update functions “copy” or “invert”. The contents of the containers will change with time. The “time” we are defining here is not the real time for the dynamics of the system. Instead, it is the time scale for a stochastic process that we use to determine the frozen core. During one time step, we remove one node from the container $``$ and determine all those nodes, to which this node is an input. A node in container $`𝒞_1`$ chooses this node as an input with probability $`1/N`$. It then becomes a frozen node. We therefore move each node of container $`𝒞_1`$ with probability $`1/N`$ into the container $``$. A node in container $`𝒞_2`$ chooses the selected frozen node as an input with probability $`2/N`$. With probability $`1/2`$, it then becomes frozen, because the frozen node is with probability $`1/2`$ in the state that fixes the output of a $`𝒞_2`$-node. If the $`𝒞_2`$-node does not become frozen, it becomes a $`𝒞_1`$-node. We therefore move each node of container $`𝒞_2`$ during the first time step with probability $`1/N`$ into the container $``$, and with probability $`1/N`$ into the container $`𝒞_1`$. Finally, a node in container $``$ chooses the selected frozen node as an input with probability $`2/N`$ and becomes a $`𝒞_1`$-node. We therefore move each node of container $``$ during the first time step with probability $`2/N`$ into the container $`𝒞_1`$. In summary, the total number of nodes, $`N`$, decreases by one during one time step, since we remove one node from container $``$, and some nodes move to a different container. The removed nodes are those frozen nodes for which we already have determined whose input they are. Then, we take the next frozen node out of container $``$ and determine its effect on the other nodes. We repeat this procedure until we cannot continue because either container $``$ is empty, or because all the other containers are empty. If container $``$ becomes empty, we are left with the nonfrozen nodes. We shall see below that most of the remaining nodes are in container $`𝒞_1`$, with the proportion of nodes left in containers $`𝒞_2`$ and $``$ vanishing in the limit $`N^{ini}\mathrm{}`$. Then, the nonfrozen nodes can be connected to a network by choosing the input(s) to every node at random from the other remaining nodes. If all containers apart from container $``$ are empty at the end, the entire network becomes frozen. This means that the dynamics of the network go to the same fixed point for all initial conditions.
Let us first describe this process by deterministic equations that neglect fluctuations around the average change of the number of nodes in the different containers. As long as all containers contain large numbers of nodes, these fluctuations are negligible, and the deterministic description is appropriate. The average change of the node numbers in the containers during one time step is
$`\mathrm{\Delta }N_r`$ $`=`$ $`{\displaystyle \frac{2N_r}{N}}`$
$`\mathrm{\Delta }N_{c_2}`$ $`=`$ $`{\displaystyle \frac{2N_{c_2}}{N}}`$
$`\mathrm{\Delta }N_{c_1}`$ $`=`$ $`{\displaystyle \frac{N_{c_1}}{N}}+{\displaystyle \frac{N_{c_2}}{N}}+{\displaystyle \frac{2N_r}{N}}`$ (1)
$`\mathrm{\Delta }N_f`$ $`=`$ $`1+{\displaystyle \frac{N_{c_1}}{N}}+{\displaystyle \frac{N_{c_2}}{N}}`$
$`\mathrm{\Delta }N`$ $`=`$ $`1`$
The number of nodes in the containers, $`N`$, can be used instead of the time variable, since it decreases by one during each step. The equation for $`N_r`$ can then be solved by going from a difference equation to a differential equation,
$$\frac{\mathrm{\Delta }N_r}{\mathrm{\Delta }N}\frac{dN_r}{dN}=\frac{2N_r}{N},$$
which has the solution
$$N_r=N^2\frac{N_r^{ini}}{(N^{ini})^2}.$$
(2)
Similarly, we find
$`N_{c_2}`$ $`=`$ $`N^2{\displaystyle \frac{N_{c_2}^{ini}}{(N^{ini})^2}}`$
$`N_f`$ $`=`$ $`N{\displaystyle \frac{N_f^{ini}N_r^{ini}}{N^{ini}}}+N^2{\displaystyle \frac{N_r^{ini}}{(N^{ini})^2}}`$
$`N_{c_1}`$ $`=`$ $`N{\displaystyle \frac{N_{c_1}^{ini}+N_{c_2}^{ini}+2N_r^{ini}}{N^{ini}}}2N^2{\displaystyle \frac{N_r^{ini}+N_{c_2}^{ini}}{(N^{ini})^2}}.`$ (3)
For $`N_f^{ini}<N_r^{ini}`$, we obtain $`N_f=0`$ at a nonzero value of $`N`$, and the number of nonfrozen nodes is proportional to $`N^{ini}`$. We are in the chaotic phase. For $`N_f^{ini}>N_r^{ini}`$, the values $`N_r`$ and $`N_{c_2}`$ will sink below 1 when $`N`$ becomes of the order $`\sqrt{N^{ini}}`$. For smaller $`N`$, there are only $``$ and $`𝒞_1`$ nodes left, and the second term contributing to $`N_f`$ and $`N_{c_1}`$ in (3) can be neglected compared to the first one. When $`N_f`$ falls below 1, there remain $`N_{c_1}=\frac{N_{c_1}^{ini}+N_{c_2}^{ini}+2N_r^{ini}}{N_f^{ini}N_r^{ini}}`$ nodes of type $`𝒞_1`$. The network is essentially frozen, with only a finite number of nonfrozen nodes in the limit $`N^{ini}\mathrm{}`$. If we now choose the inputs for these nodes, we obtain simple loops with trees rooted in the loops. This property of the frozen phase was also found in Socolar and Kauffman (2003).
For the critical networks that this paper focuses on, we have $`N_f^{ini}=N_r^{ini}=\beta N^{ini}`$, and the stochastic process stops at $`N_f=1=\beta N^2/N^{ini}`$. This means that
$$N^{end}=\sqrt{\frac{N^{ini}}{\beta }}.$$
(4)
The number of nonfrozen nodes would scale with the square root of the network size if the deterministic approximation to the stochastic process was exact. We shall see below that including fluctuations changes the exponent from $`1/2`$ to $`2/3`$. The final number of $`𝒞_2`$-nodes for the deterministic process for the critical networks is $`\gamma /\beta `$, which is independent of network size, and the final number of $``$-nodes vanishes due to $`N_r=N_f`$. We shall see below that the fluctuations change these two results to a $`(N^{ini})^{1/3}`$-dependence.
Introducing $`n=N/N^{ini}`$ and $`n_j=N_j/N^{ini}`$ for $`j=r,f,c_1,c_2`$, equations (3) simplify to (using $`N_r^{ini}=N_f^{ini}`$)
$`n_r`$ $`=`$ $`\beta n^2=n_f`$
$`n_{c_2}`$ $`=`$ $`\gamma n^2`$
$`n_{c_1}`$ $`=`$ $`n2\beta n^2\gamma n^2.`$
This means that our stochastic process remains invariant (in the deterministic approximation) when the initial number of nodes in the containers and the time unit are all multiplied by the same factor. For small $`n`$, the majority of nodes are in container $`𝒞_1`$, since $`n_{c_1}=n𝒪(n^2)`$. Now, if we choose a sufficiently large $`N^{ini}`$, $`n`$ reaches any given small value while $`N_f=N_r=\beta n^2N^{ini}`$ is still large enough for a deterministic description. We can therefore assume that for sufficiently large networks $`N_f/N=\beta n`$ becomes small before the effect of the noise becomes important. This assumption will simplify our calculations below.
## 4 The effect of fluctuations
The number of nodes in container $`𝒞_1`$ that choose a given frozen node as an input is Poisson distributed with a mean $`N_{c_1}/N`$ and a variance $`N_{c_1}/N`$. We now assume that $`n`$ is small at the moment where noise becomes important, i.e., that the variance of the noise $`N_{c_1}/N=n_{c_1}/n=1(2\beta +\gamma )n=1𝒪(n)`$ is unity. The number of nodes in containers $`𝒞_2`$ and $``$ that choose a given frozen node as an input is Poisson distributed with a mean and a variance $`2(N_{c_2}+N_r)/N`$. The fluctuation around the mean can be neglected as this noise term is very small compared to $`N_r`$ and $`N_{c_2}`$, the final values of which are large for sufficiently large $`N^{ini}`$. We therefore obtain the stochastic version of equations (1)
$`\mathrm{\Delta }N_r`$ $`=`$ $`{\displaystyle \frac{2N_r}{N}}`$
$`\mathrm{\Delta }N_{c_2}`$ $`=`$ $`{\displaystyle \frac{2N_{c_2}}{N}}`$
$`\mathrm{\Delta }N_f`$ $`=`$ $`{\displaystyle \frac{N_r}{N}}{\displaystyle \frac{N_f}{N}}+\xi `$
$`\mathrm{\Delta }N`$ $`=`$ $`1`$ (5)
The random variable $`\xi `$ has zero mean and unit variance. As long as the $`n_j`$ change little during one time step, we can summarize a large number $`T`$ of time steps into one effective time step, with the noise becoming Gaussian distributed with zero mean and variance $`T`$. Exactly the same process would result if we summarized $`T`$ time steps of a process with Gaussian noise of unit variance. For this reason, we can choose the random variable $`\xi `$ to be Gaussian distributed with unit variance.
Compared to the deterministic case, the equations for $`N_r`$ and $`N_{c_2}`$ are unchanged, and we have again $`N_r=N^2N_r^{ini}/(N^{ini})^2`$ and $`N_{c_2}=N^2N_{c_2}^{ini}/(N^{ini})^2`$. Inserting the solution for $`N_r`$ into the equation for $`N_f`$, we obtain
$$\frac{dN_f}{dN}=\frac{N_f}{N}+\frac{\beta N}{N^{ini}}+\xi $$
(6)
with the step size $`dN=1`$ and $`\xi ^2=1`$. (In the continuum limit $`dN0`$ the noise correlation becomes $`\xi (N)\xi (N^{})=\delta (NN^{})`$). This is a Langevin-equation, and we will now derive the corresponding Fokker-Planck-equation. Let $`P(N_f,N)`$ be the probability that there are $`N_f`$ nodes in container $``$ at the moment where there are $`N`$ nodes in total in the containers. This probability depends on the initial node number $`N_{ini}`$, and on the parameter $`\beta `$. The sum
$$\underset{N_f=1}{\overset{\mathrm{}}{}}P(N_f,N)_0^{\mathrm{}}P(N_f,N)𝑑N_f$$
is the probability that the stochastic process is not yet finished, i.e. the probability that $`N_f`$ has not yet reached the value 0 at the moment where the total number of nodes in the containers has decreased to the value $`N`$. Since systems that have reached $`N_f=0`$ are removed from the ensemble, we have to impose the absorbing boundary condition $`P(0,N)=0`$. Let $`g(\mathrm{\Delta }N_f|N_f,N)`$ denote the probability that $`N_f`$ decreases by $`\mathrm{\Delta }N_f`$ during the next step, given the values of $`N_f`$ and $`N`$.
We have
$`P(N_f,N1)=`$
$`{\displaystyle _0^{\mathrm{}}}P(N_f+\mathrm{\Delta }N_f,N)g(\mathrm{\Delta }N_f|N_f+\mathrm{\Delta }N_f,N)d(\mathrm{\Delta }N_f)`$
$`={\displaystyle _0^{\mathrm{}}}[P(N_f,N)g(\mathrm{\Delta }N_f|N_f,N)`$
$`+{\displaystyle \frac{}{N_f}}(P(N_f,N)g(\mathrm{\Delta }N_f|N_f,N))\mathrm{\Delta }N_f`$
$`+{\displaystyle \frac{^2}{2^2N_f^2}}(P(N_f,N)g(\mathrm{\Delta }N_f|N_f,N))(\mathrm{\Delta }N_f)^2`$
$`+\mathrm{}]d(\mathrm{\Delta }N_f)`$
$`=P(N_f,N)+{\displaystyle \frac{}{N_f}}(P(N_f,N)\mathrm{\Delta }N_f)+`$
$`{\displaystyle \frac{^2}{2N_f^2}}(P(N_f,N)(\mathrm{\Delta }N_f)^2)+\mathrm{}`$
The mean change $`\mathrm{\Delta }N_f`$ during one step is $`\mathrm{\Delta }N_f=\frac{N_f}{N}+\frac{\beta N}{N^{ini}}`$, and the mean square change is $`(\mathrm{\Delta }N_f)^21`$.
This gives the Fokker-Planck equation for our stochastic process
$$\frac{P}{N}=\frac{}{N_f}\left(\frac{N_f}{N}+\frac{\beta N}{N^{ini}}\right)P+\frac{1}{2}\frac{^2P}{N_f^2}.$$
(7)
We introduce the variables
$$x=\frac{N_f}{\sqrt{N}}\text{ and }y=\frac{N}{(N^{ini}/\beta )^{2/3}}$$
(8)
and the function $`f(x,y)=(N^{ini}/\beta )^{1/3}P(N_f,N)`$. We will see below that $`f(x,y)`$ does not depend explicitely on the parameters $`N^{ini}`$ and $`\beta `$ with this definition. The Fokker-Planck equation then becomes
$$y\frac{f}{y}+f+\left(\frac{x}{2}+y^{3/2}\right)\frac{f}{x}+\frac{1}{2}\frac{^2f}{x^2}=0.$$
(9)
Let $`W(N)`$ denote the probability that $`N`$ nodes are left at the moment where $`N_f`$ reaches the value zero. It is
$`W(N)`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}P(N_f,N)𝑑N_f{\displaystyle _0^{\mathrm{}}}P(N_f,N1)𝑑N_f.`$
Consequently,
$`W(N)`$ $`=`$ $`{\displaystyle \frac{}{N}}{\displaystyle _0^{\mathrm{}}}P(N_f,N)𝑑N_f`$ (10)
$`=`$ $`(N^{ini}/\beta )^{1/3}{\displaystyle \frac{}{N}}\sqrt{N}{\displaystyle _0^{\mathrm{}}}f(x,y)𝑑x`$
$`=`$ $`(N^{ini}/\beta )^{2/3}{\displaystyle \frac{}{y}}\sqrt{y}{\displaystyle _0^{\mathrm{}}}f(x,y)𝑑x`$
$``$ $`(N^{ini}/\beta )^{2/3}G(y)`$
with a scaling function $`G(y)`$. $`W(N)`$ must be a normalized function, $`_0^{\mathrm{}}W(N)𝑑N=_0^{\mathrm{}}G(y)𝑑y=1`$. This condition is independent of the parameters of the model, and therefore $`G(y)`$ and $`f(x,y)`$ are independent of them, too, which justifies our choice of the prefactor in the definition of $`f(x,y)`$. By integrating equation (9) over $`x`$ from 0 to infinity and by using $`f(0,y)=f(\mathrm{},y)=0`$ we obtain
$$\sqrt{y}\frac{}{y}\sqrt{y}_0^{\mathrm{}}f𝑑x\frac{1}{2}\frac{f}{x}|_{x=0}=0,$$
which gives us a second relation between $`f(x,y)`$ and $`G(y)`$:
$$\sqrt{y}G(y)=\frac{1}{2}\frac{f}{x}|_{x=0}.$$
(11)
The mean number of nonfrozen nodes is
$$\overline{N}=_0^{\mathrm{}}NW(N)𝑑N=(N^{ini}/\beta )^{2/3}_0^{\mathrm{}}G(y)y𝑑y,$$
(12)
which is proportional to $`(N^{ini}/\beta )^{2/3}`$. We did not succeed in extracting an explicit expression for the function $`G(y)`$. It can be determined by running the stochastic process described by the equations (5) on the computer. The result is shown in Figure 1, and an almost perfect fit to this result is given by
$$G(y)0.25e^{y^3/2}(10.5\sqrt{y}+3y)/\sqrt{y}.$$
(13)
For small $`y`$, the data show a power law $`G(y)y^{1/2}`$. We can obtain this power law analytically by solving the Fokker-Planck equation (9) in the limit of small $`y`$. In this limit, the term proportional to $`y^{3/2}`$ can be dropped, and we have the simpler equation
$$y\frac{f}{y}+f+\frac{x}{2}\frac{f}{x}+\frac{1}{2}\frac{^2f}{x^2}=0.$$
(14)
The general solution has the form $`f(x,y)=_\nu c_\nu y^\nu f_\nu (x)`$, with the functions $`f_\nu `$ satisfying
$$2(\nu +1)f_\nu +xf_\nu ^{}+f_\nu ^{\prime \prime }=0.$$
(15)
The solution is
$$e^{\frac{x^2}{2}}f_\nu (x)=C_1H_{1+2\nu }\left(\frac{x}{\sqrt{2}}\right)+C_2{}_{1}{}^{}F_{1}^{}(\nu \frac{1}{2};\frac{1}{2};\frac{x^2}{2})$$
with two constants $`C_1`$ and $`C_2`$, and with $`H`$ denoting the Hermitian functions, and $`{}_{1}{}^{}F_{1}^{}`$ the appropriate hypergeometric functions. We expect $`f`$ to be analytical in $`y`$ for small $`y`$, which means that $`\nu =0,1,2,\mathrm{}`$. For sufficiently small $`y`$, only the term $`\nu =0`$ contributes, and due to the absorbing boundary condition we have $`C_2=0`$. We obtain therefore for small $`y`$
$$f(x,y)=c_0xe^{x^2/2}.$$
(16)
From our numerical result (13), together with (11), we find $`c_0=0.5`$. Inserting Eq. (16) into Eq. (10), we obtain for small $`N`$
$$W(N)=\left(\frac{N^{ini}}{\beta }\right)^{1/3}\frac{c_0}{2\sqrt{N}}.$$
(17)
In Eq. (16), the function $`f(x,y)`$ is independent of $`y`$. This means that for sufficiently small $`N`$ the function $`P(N_f,N)`$ depends only on the ratio $`N_f/\sqrt{N}`$. This is also confirmed by our computer simulations (see Fig. 2).
We can obtain a set of solutions of Eq. (9) with the Ansatz $`f(x,y)=_\nu y^\nu \stackrel{~}{f}_\nu (z)`$ with $`z=xy^{3/2}`$. The resulting equation for $`\stackrel{~}{f}_\nu `$, is identical to Eq. (15) for $`f_\nu `$, which was valid for small $`y`$. However, an analytical expression for the function $`G(y)`$ can only be given if an expansion of the initial condition $`P(N_f,N)=\delta (N_f\beta N^{ini})`$ in terms of known solutions can be found.
The probability $`W_r(N_r)`$ that $`N_r`$ nodes are left in container $``$ at the moment where container $``$ becomes empty, is obtained from the relation
$$N_r=N^2N_r^{ini}/(N^{ini})^2.$$
Defining
$$s=\frac{N_r}{(N^{ini}/\beta )^{1/3}}=y^2$$
and
$$F(s)=\frac{G(\sqrt{s})}{2\sqrt{s}},$$
(18)
and remembering $`W(N)dN=W_r(N_r)dN_r`$, we find
$$W_r(N_r)=(N^{ini}/\beta )^{1/3}F(s).$$
(19)
The mean number of nodes left in in container $``$ is
$`\overline{N}_r`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}W_r(N_r)N_r𝑑N_r=(N^{ini}/\beta )^{1/3}{\displaystyle _0^{\mathrm{}}}sF(s)𝑑s`$ (20)
$`=`$ $`(N^{ini}/\beta )^{1/3}{\displaystyle _0^{\mathrm{}}}y^2G(y)𝑑y.`$
The number of nodes left in container $`𝒞_2`$ is $`N_{c_2}=(\gamma /\beta )N_r`$.
We thus have shown that the number of nonfrozen nodes scales with network size $`N^{ini}`$ as $`(N^{ini})^{2/3}`$, with most of these nodes receiving only one input from other nonfrozen nodes. The number of nonfrozen nodes receiving two inputs from nonfrozen nodes scales as $`(N^{ini})^{1/3}`$. We have found scaling functions that describe the probability distribution for these two types of nodes in the limit of large network size. Our next task will be to connect these nonfrozen nodes to a network. This is a reduced network, where all frozen nodes have been cut off.
## 5 Relevant nodes
Let us start from the result obtained from the stochastic process of the previous two sections. Each time we run this process we obtain $`N`$ nonfrozen nodes. Out of these, $`N_r`$ ($`N_{c_2}`$) nodes receive input from two other nonfrozen nodes and have a reversible (canalizing $`𝒞_2`$) update function. We define the parameter
$$a=\frac{N_r+N_{c_2}}{\sqrt{N}}=(1+\gamma /\beta )y^{3/2},$$
(21)
which has a probability distribution $`f(a)`$ that is determined from the condition $`f(a)da=G(y)dy`$,
$$f(a)=\frac{2}{3a^{1/3}(1+\gamma /\beta )^{2/3}}G\left(\left(\frac{a}{1+\gamma /\beta }\right)^{2/3}\right).$$
(22)
Just as $`G(y)`$, the function $`f(a)`$ is the exact probability distribution only in the thermodynamic limit $`N^{ini}\mathrm{}`$. We determine the relevant nodes by a stochastic process that removes iteratively nodes that are not relevant. Each of the $`N`$ nonfrozen nodes chooses its input(s) at random from the nonfrozen nodes. There are altogether $`N(1+a/\sqrt{N})`$ inputs to be chosen, and consequently the nonfrozen nodes have together $`N(1+a/\sqrt{N})`$ outputs. The number of outputs of a node is Poisson distributed with the mean value $`(1+a/\sqrt{N})`$. The fraction $`\mathrm{exp}(1a/\sqrt{N})`$ of nodes have no output. They are the leaves of the trees of the network of nonfrozen nodes, and we therefore know that they are not relevant. We put them in container number 1. Their number will change during the stochastic process that determines the relevant nodes. The other nodes are placed in container number 2. Their number is $`N_l`$ (“labelled”), and it will be reduced until only the relevant nodes are left. The total number of outputs of the nodes in container 2 is initially $`N(1+a/\sqrt{N})`$, while their total number of inputs is $`N(1+a/\sqrt{N})(1\mathrm{exp}(1a/\sqrt{N}))`$. Now, we remove one node from container 1 and connect its input(s) at random to the outputs of the nodes in container 2. The chosen output(s) are cut off. If a node whose output is cut off has no other output left, we move the node from container 2 to container 1. It cannot be a relevant node since relevant nodes influence other relevant nodes. We iterate this procedure, until there is no node left in container 1. The nodes remaining in container 2 are the relevant nodes. During the entire process, the number of outputs in container 2 is identical to the number of inputs in container 1 and 2. As long as container 1 is not empty, there are more outputs in container 2 than inputs, and only when the process is finished do the two numbers become identical. We can therefore simplify the stochastic process by removing container 1 altogether. We simply have to continue cutting of outputs from nodes in container 2 and removing nodes with no outputs, until the total number of outputs of the nodes in container 2 has become identical to their total number of inputs. The remaining nodes are relevant, and we have then $`N_l^{final}N_{rel}`$. These nodes can then be connected to a network by connecting the inputs and outputs pairwise.
In order to derive analytical results, it is useful to run this process backwards. Starting with $`N`$ nodes with no outputs, adding outputs at random will eventually generate the Poisson distribution of the number of outputs per node that we have started with. The reverse stochastic process is therefore defined by the following rule: Begin with an empty container (former container 2) and $`N`$ nodes outside the container. Most of these nodes have one input, and the fraction $`a/\sqrt{N}`$ have two inputs. Add an output to a randomly chosen node. Put this node in the container. Add another output to a randomly chosen node (choosing every node with equal probability, whether the node is inside or outside the container). If a node from outside the container is chosen, put it in the container. Eventually, the total number of outputs in the container will become larger than the total number of inputs in the container. The container contains the relevant nodes at the moment when the inputs equal the outputs for the last time.
In order to show that the number of relevant nodes scales with $`\sqrt{N}`$, we define a scaling variable
$$t=\frac{N_l}{\sqrt{N}}.$$
During one step, an output is added to nodes that are already in the container with probability $`N_l/N`$. Let $`N_o`$ count the number of outputs that have been added to nodes in the container, i.e., $`N_o=`$(total number of outputs in the container) $`N_l`$. Then the average rate of increase of $`N_o`$ is given for sufficiently large $`N`$ by
$$\frac{dN_o}{dN_l}=\frac{N_l}{N},$$
or
$$\frac{dN_o}{dt}=t.$$
Let $`N_i`$ count the number of nodes in the container with two inputs. Their rate of increase is
$$\frac{dN_i}{dN_l}=\frac{a}{\sqrt{N}},$$
or
$$\frac{dN_i}{dt}=a.$$
Consequently, the probability distribution for $`N_o`$ is given by
$$P_o(N_o|t)=\frac{1}{N_o!}e^{t^2/2}\left(\frac{t^2}{2}\right)^{N_o},$$
(23)
and the probability distribution for $`N_i`$ is given by
$$P_i(N_i|t)=\frac{1}{N_i!}e^{at}\left(at\right)^{N_i}.$$
(24)
The stochastic process can be viewed as a random walk that steps to the right with a rate $`t`$ and to the left with a rate $`a`$. It is finished when $`N_i=N_o`$ for the last time, i.e. when the walk leaves the origin for the last time. We determined the probability distribution $`𝒞_a(t)`$ for this last exit time from the origin by a computer simulation. It is shown in Fig. 4 for $`a=1`$. For small $`t`$, it increases linearly in $`t`$, because the probability of making a step to the right is proportional to $`t`$ for small times.
For $`a=0`$, we can obtain an analytical result from the relation
$$𝒞_0(t)=\frac{P_o(0,t)}{t}=te^{t^2/2}.$$
(25)
Since we were able to write the stochastic process in terms of $`t`$ and $`a`$ alone, the probability distribution for the number of relevant nodes depends only on the combination $`N_{rel}/\sqrt{N}`$ and on the parameter $`a`$,
$$p_a(N_{rel})dN_{rel}=𝒞_a\left(N_{rel}/\sqrt{N}\right)dN_{rel}/\sqrt{N}.$$
(26)
The relation between $`N`$ and $`a`$ is obtained using Eq. (8) and (21):
$$\sqrt{N}=a^{1/3}\left(\frac{N^{ini}}{\beta +\gamma }\right)^{1/3}.$$
Taking into account the probability distribution (22) of the parameter $`a`$, we obtain the scaling behavior of the number of relevant nodes,
$$p(N_{rel})=_0^{\mathrm{}}𝑑af(a)𝒞_a\left(\frac{N_{rel}a^{1/3}}{(N^{ini}/(\beta +\gamma ))^{1/3}}\right)\left(\frac{\beta +\gamma }{aN^{ini}}\right)^{1/3}.$$
(27)
The error made by taking the upper limit of the integral to infinity vanishes for $`N^{ini}\mathrm{}`$. We introduce the scaling variable
$$z=\frac{N_{rel}}{\left(\frac{N^{ini}}{\beta +\gamma }\right)^{1/3}},$$
(28)
which has then the following probability distribution
$$P(z)=_0^{\mathrm{}}𝑑a\frac{f(a)}{a^{1/3}}𝒞_a\left(\frac{z}{a^{1/3}}\right).$$
(29)
The probability distribution for the number of relevant nodes depends for large $`N^{ini}`$ only on the scaling variable $`z`$. We determined numerically the function $`P(z)`$ by combining the two stochastic processes described in this paper. First, we determined a value of $`a`$ using the process of Section 4. Then, we used this value of $`a`$ to determine the last exit time of the stochastic process of this section, giving a value of $`z`$. The shape of the curves $`P(z)`$ depends on the value of $`\gamma /\beta `$, and the results are shown in Fig. 5 for $`\gamma /\beta =0`$ and $`\gamma /\beta =4`$, which is the original Kauffman model, where each update function has the same weight. It is easy to check analytically that
$$\underset{z0}{lim}P(z)=\sqrt{2\pi }/4(1+\gamma /\beta )^{1/3}.$$
The mean number of relevant nodes is
$$\overline{N}_{rel}=_0^{\mathrm{}}N_{rel}p(N_{rel})𝑑N_{rel}=\left(\frac{N^{ini}}{\beta +\gamma }\right)^{1/3}_0^{\mathrm{}}zP(z)𝑑z,$$
(30)
i.e., it is proportional to $`(N^{ini})^{1/3}`$. Finally, let us give the probability distribution for the number of relevant nodes with two relevant inputs. Let $`m`$ denote the number of relevant nodes with two relevant inputs and $`\stackrel{~}{P}(m;z)dz`$ the probability of having the number of relevant nodes in the interval $`[N_{rel}(z),N_{rel}(z+dz)]`$, with $`m`$ of them having two relevant inputs. Using Equations (23) and (24), we can express $`\stackrel{~}{P}`$ as
$`\stackrel{~}{P}(m;z)`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}𝑑a{\displaystyle \frac{f(a)}{a^{1/3}}}𝒞_a\left({\displaystyle \frac{z}{a^{1/3}}}\right)`$
$`\times {\displaystyle \frac{P_o\left(m|za^{1/3}\right)P_i\left(m|za^{1/3}\right)}{_lP_o\left(l|za^{1/3}\right)P_i\left(l|za^{1/3}\right)}}.`$
As $`P_o`$ and $`P_i`$ decay exponentially fast with increasing $`m`$, the mean number of relevant nodes with two inputs is finite.
## 6 Conclusions
In this paper, we have obtained the asymptotic probability distributions in the limit of large network size for the number of nonfrozen nodes, the number of nonfrozen nodes with two nonfrozen inputs, the number of relevant nodes, and the number of relevant nodes with two relevant inputs. The mean values of these quantities scale with network size $`N^{ini}`$ as a power law in $`N^{ini}`$, with the exponent being $`2/3`$, $`1/3`$, $`1/3`$, and $`0`$ respectively. The implications of the results are manifold.
First, the notion that these networks are “critical” is now corroborated by the existence of power laws and scaling functions. Originally, it was expected that the quantities that display the scaling behavior should be the attractors of the network Kauffman (1969). In the meantime, it has become clear that mean attractor numbers do not obey power laws Samuelsson and Troein (2003). It is the number of nonfrozen and relevant nodes that show scaling behavior.
Next, let us compare the results to those of critical $`K=1`$ networks. A $`K=1`$ critical network with $`N`$ nodes corresponds to the nonfrozen part of a critical $`K=2`$ network for $`a=0`$. In this case, the probability distribution of the number of relevant nodes is given by Eq. (26) with $`a=0`$,
$$p_0(N_{rel})=\frac{1}{\sqrt{N}}𝒞_0\left(\frac{N_{rel}}{\sqrt{N}}\right)=\frac{N_{rel}}{N}e^{N_{rel}^2/2N}.$$
(32)
The mean number of relevant nodes is proportional to $`\sqrt{N}`$. When these relevant nodes are connected to a network by pairwise connecting the inputs and outputs, one obtains a set of simple loops. From Drossel et al. (2005), we know that there is a mean number of $`\mathrm{ln}\sqrt{N}`$ loops and that the number of loops of length $`l`$ in a critical $`K=1`$ network is Poisson distributed with a mean $`1/l`$ for $`l\sqrt{N}`$. This can be easily explained by consindering the process of connecting inputs and outputs: We begin with a given node and draw the node that provides its input from all possible nodes. Then, we draw the node that provides the input to the newly chosen node, etc., until the first node is chosen and a loop is formed. For small loop size, the probability that the loop is closed after the addition of the $`l`$th node is $`1/N_{rel}`$. Therefore, the probability that a given node is on a loop of size $`l`$ is $`1/N_{rel}`$, and the mean number of nodes on loops of size $`l`$ is 1, and the number of loops of length $`l`$ is Poisson distributed with a mean $`1/l`$ for sufficiently small $`l`$.
Now, the $`K=2`$ critical networks have of the order of $`(N^{ini})^{1/3}`$ relevant nodes, with only a finite number of them having two relevant inputs. The relevant components are constructed from the relevant nodes by pairwise connecting inputs and outputs. In the asymptotic limit of very large $`N^{ini}`$ that we are considering, the probability that a randomly chosen relevant node has two inputs or two outputs vanishes. Let us again construct a component by starting with one node and choosing its input node etc., until the component is finished. If the component is small, it consists almost certainly only of nodes with one input and one output and is therefore a simple loop. There is no difference between the statistics of the small relevant components of a $`K=1`$ critical network, and the number of loops of length $`l`$ is Poisson distributed with a mean $`1/l`$. The total number of relevant nodes in loops of size $`ll_c`$ with $`l_c=ϵ(N^{ini})^{1/3}`$ (with a small $`ϵ`$) is $`l_c`$, and it is a small proportion of all nodes. If there were no nodes with two inputs or outputs, the number of components larger than $`l_c`$ would be $`(\mathrm{ln}N_{rel}\mathrm{ln}l_c)=\mathrm{ln}(1/ϵ)`$. The additional links may reduce this number, which is in any case finite. Since these large components contain almost all nodes, they contain almost all relevant nodes with two inputs or outputs.
From these findings, we can obtain results for the attractors of $`K=2`$ critical networks. The numbers and lengths of attractors are determined by the relevant components. We now argue that the mean number and length of attractors increases faster than any power law. If we remove the components of size larger than $`l_c`$ and determine the mean number and length of attractors for this reduced relevant network, we have a lower bound to the correct numbers. Now, the reduced relevant network of a $`K=2`$ system is identical to that of a critical $`K=1`$ system (where the critical loop size is $`l_c=ϵ\sqrt{N}`$). In Drossel et al. (2005), it was proven that the mean number and length of attractors for such a reduced $`K=1`$ system increases faster than any power law with network size. We therefore conclude that the same is true for critical $`K=2`$ networks.
Earlier, Samuelsson and Troein Samuelsson and Troein (2003) have derived analytically an exact expression for the number of attractors of length $`L`$ of a critical $`K=2`$ network in the limit of large $`N^{ini}`$, and they have pointed out that this implies that the mean number of attractors increases faster than any power law with $`N^{ini}`$. Using their calculation, it has recently been shown Drossel (2005) that there is a close relationship between $`K=1`$ critical networks and the nonfrozen part of $`K=2`$ critical networks, and that the results of Samuelsson and Troein (2003) can be most naturally interpreted if the relevant components of these two networks look identical for component sizes below the above-given cutoffs. This interpretation is placed on a firm foundation by the present paper. |
warning/0506/math0506135.html | ar5iv | text | # On differentiable compactifications of the hyperbolic space
## Introduction
### 0.1 Goal
The group $`G=\text{SO}\text{0}\text{(n,1)}`$ acts on the $`n`$-dimensional open ball as the isometries for the hyperbolic riemanniann metric (we denote this action by $`\mathrm{𝗂𝗌𝗈𝗆}`$). We study the differentiable compactifications of this action into an action of $`G`$ on the closed ball $`\overline{𝔹}^n`$, that is to say the differentiable actions of $`G`$ on the closed ball such that the restriction to the open ball is differentiably conjugate to $`\mathrm{𝗂𝗌𝗈𝗆}`$.
We concentrate on two levels of regularity: “differentiable” shall always mean smooth or real analytic.
One analytic compactification is given by the continuous extension of the action of $`G`$ in the Klein model of the hyperbolic space. We shall give later on more details about this compactification, called “projective” and denoted by $`\mathrm{𝗉𝗋𝗈𝗃}`$.
It appears that up to an analytic change of coordinates, there is a countable familly of analytic compactifications. The classification is more complicated in the smooth case, but we get a simple description of all smooth compactifications.
To describe the compactifications, we use half-space charts of $`\overline{𝔹}^n`$ missing only one point.
###### Theorem 1
In some half space chart the conjugation of $`\mathrm{𝗉𝗋𝗈𝗃}`$ by a change of coordinates of the form:
$$\phi _p:(x_1,\mathrm{},x_{n1},y)(x_1,\mathrm{},x_{n1},y^p)$$
(1)
(where $`p`$ is a positive integer) is well defined in the open ball and continuously extendable into an analytic action on the closed ball. Thus for each positive integer $`p`$ we get an analytic compactification of $`\mathrm{𝗂𝗌𝗈𝗆}`$. Moreover, any analytic compactification is analytically conjugate to one of these and no two different $`p`$’s give conjugate compactifications.
We shall prove this result when $`n3`$ ; the case $`n=2`$ can be deduced from works of Schneider and Stowe exposed by Mitsumatsu .
###### Theorem 2
In some half space chart the conjugation of $`\mathrm{𝗉𝗋𝗈𝗃}`$ by a change of coordinates of the form:
$$\phi _f:(x_1,\mathrm{},x_{n1},y)(x_1,\mathrm{},x_{n1},f(y))$$
(2)
(where $`f:^+^+`$ is a homeomorphism that is smooth and a diffeomorphism of $`^+^{}`$) is well defined in the open ball and is continuously extendable into a smooth action on the closed ball if and only if
$$f/f^{}\text{ is smooth.}$$
(3)
Thus for each map $`f`$ satisfying (3) we get a smooth compactification of $`\mathrm{𝗂𝗌𝗈𝗆}`$. Moreover, if $`n3`$ any smooth compactification is analytically conjugate to one of these. Two different $`f`$’s give conjugate compactifications if and only if they are smoothly conjugate.
Condition (3) is automatically satisfied if $`f`$ is non-flat (i.e. has some non-zero derivative in 0).
Remarks
1. By “$`f`$ is smooth” we mean that it can be prolonged into a smooth map defined on $`]\epsilon ,+\mathrm{}[`$ where $`\epsilon `$ is some positive real number. Equivalently, $`f`$ is smooth on $`^+`$ if and only if all its derivatives converge in $`0`$,
2. some flat maps satisfy condition (3) while others do not. For example, if $`f_1=x\mathrm{exp}(x^2)`$ and $`f_2=x\mathrm{exp}(x^{\frac{3}{2}})`$ we have $`(f_1/f_1^{})(x)=\frac{1}{2}x^3`$ but $`(f_2/f_2^{})(x)=\frac{2}{3}x^{\frac{5}{2}}`$,
3. in general, $`f`$ is singular at $`0`$. Otherwise, the compactification given by $`\phi _f`$ is conjugate to $`\mathrm{𝗉𝗋𝗈𝗃}`$,
4. in the two-dimensional case, it appears that there are only two smooth compactifications of $`\mathrm{𝗂𝗌𝗈𝗆}`$ (namely the projective one and the conformal one, see 0.2.1 for a definition) that are algebraic, i.e. given by a linear representation of $`G`$ after projectivization and restriction to an embedded manifold that is a union of orbits. Thanks to Theorem 2 it is easy to generalize this result to higher dimensions (an algebraic action of SO<sub>0</sub>(1,n) on the closed $`n`$-ball gives by restriction to a totally geodesic plane an algebraic action of $`\text{SL}_2\text{(}\text{)}`$ on the closed $`2`$-ball).
### 0.2 Recalls and notations
#### 0.2.1 The projective and conformal compactifications
Two models of the isometric action of SO<sub>0</sub>(n,1) on the hyperbolic space are well-known: the Poincaré ball and the Klein ball.
Let us recall the construction of these models. Put on $`^{n+1}`$ a Lorentzian quadratic form $`Q=x_1^2+x_2^2+\mathrm{}+x_n^2y^2`$ and consider the hypersurface $`H`$ defined by $`Q=1`$ and $`y>0`$. $`Q`$ induces a scalar product on each tangent space of $`H`$, thus defines a riemannian structure: the hyperbolic space $`^n`$, on which SO<sub>0</sub>(n,1) acts naturally by isometries. A central projection of center $`0`$ of $`H`$ on the disc of radius 1 of the hyperpane $`y=1`$ gives the Klein ball model. If we project vertically the Klein ball on the upper half-sphere of radius 1 of $`^{n+1}`$ and then project the half-sphere stereographically from the point $`(0,0,\mathrm{},0,1)`$ on the plane $`y=0`$, we get the Poincaré ball model (see figure 1).
In both of these models, the isometric action of SO<sub>0</sub>(n,1) appears to admit continuous extensions to the closed ball into an analytic action.
The extension in the Poincaré model is called the conformal action: it preserves the conformal euclidian structure of the closed unit ball. We denote this action by $`\mathrm{𝖼𝗈𝗇𝖿}`$, and the action of an element $`gG`$ by $`\mathrm{𝖼𝗈𝗇𝖿}_g`$.
The extension in the Klein model is called the projective action: it preserves the projective structure of the closed ball, subset of $`(^{n+1})`$ (in particular, the projective action maps straight lines into straight lines). We denote this action by $`\mathrm{𝗉𝗋𝗈𝗃}`$, and the action of an element $`gG`$ by $`\mathrm{𝗉𝗋𝗈𝗃}_g`$. We also denote by $`\mathrm{𝗉𝗋𝗈𝗃}_X`$ the vector field given by the projective action of an element $`X`$ of the Lie algebra $`𝔤`$ of $`G`$.
It is easy to see that the conformal and projective actions are topologically conjugate, but they are not differentiably conjugate.
Let us give a purely geometrical proof.
In the conformal model, two asymptotic geodesics are always tangent but in the projective model, for any direction in the tangent space of a point of the boundary there is exactly one geodesic tangent to it (we shall see that this property of $`\mathrm{𝗉𝗋𝗈𝗃}`$ plays a fundamental role in Theorems 1 and 2).
Therefore, the group of the parabolic elements of $`G`$ which stabilize a given point of the boundary has a common proper direction transversal to the boundary (in particular their linear parts are simultaneously diagonalisable) for the conformal action but not for the projective one, hence these two compactifications cannot be differentiably conjugate.
#### 0.2.2 Half-space charts
We denote by $`^{n+}`$ the open $`n`$-dimensional half space and by $`\overline{}^{n+}`$ its closure in $`^n`$. We shall use the canonical coordinates system $`(x_1,x_2,\mathrm{},x_{n1},y)`$, therefore $`^{n+}=\left\{(x_1,\mathrm{},x_{n1},y)^n;y>0\right\}`$.
The extension to the boundary of the isometric action of the Poincaré half-space is the conformal action written in some half-space chart of the closed ball. We denote this chart by $`𝒫𝒞`$, it is given by: a central projection of the Poincaré disc from the south pole $`(0,0,\mathrm{},0,1)`$ on the upper half-sphere, composed with a stereographic projection of center $`(0,0,\mathrm{},0,1,0)`$ on the vertical $`n`$-dimensional half-space defined by $`x_n=0`$ (identified with $`^{n+}`$).
If we compose the vertical projection of the Klein Ball on the upper half-sphere with the stereographic projection of center $`(0,0,\mathrm{},0,1,0)`$ on the vertical $`n`$-dimensional half-space defined by $`x_n=0`$ and then with the following map:
$$\begin{array}{cccc}\hfill \phi _2:& ^{n+}& & ^{n+}\\ & (x_1,\mathrm{},x_{n1},y)& & (x_1,\mathrm{},x_{n1},y^2)\end{array}$$
it appears that we get a chart, denoted by $`𝒦𝒞`$. An explicit change of coordinates defining this chart from the projective coordinates is given by:
$$[1:x_1:x_2:\mathrm{}:x_{n1}:y](\frac{x_1}{1y},\frac{x_2}{1y},\mathrm{},\frac{1{\displaystyle \underset{i=1}{\overset{n1}{}}}x_i^2y^2}{(1y)^2})$$
(4)
Hence $`\mathrm{𝖼𝗈𝗇𝖿}`$ corresponds in Theorem 1 to the case $`p=2`$. In particular $`\mathrm{𝖼𝗈𝗇𝖿}`$ and $`\mathrm{𝗉𝗋𝗈𝗃}`$ are Hölder-$`\frac{1}{2}`$ conjugate.
## 1 Topological uniqueness
We saw that the conformal and projective actions are topologically conjugate. More generally we prove the following fact:
###### Proposition 3
There is up to topological conjugacy only one continuous action of $`G`$ on the closed ball such that the action on the interior is homeomorphic to $`\mathrm{𝗂𝗌𝗈𝗆}`$.
This uniqueness enables us to use the following definition:
###### Definition 4
Throughout this paper, by a $`𝒞^k`$ compactification ($`k`$ is $`\mathrm{}`$ or $`\omega `$) of $`\mathrm{𝗂𝗌𝗈𝗆}`$ we shall mean a homeomorphism $`\phi `$ of the closed ball which is a $`𝒞^k`$ diffeomorphism in the interior and such that for each $`gG`$, $`\phi ^1\mathrm{𝗉𝗋𝗈𝗃}_g\phi `$ (defined on the open ball) admits a $`𝒞^k`$ extension to the closed ball.
In other words, we do not distinguish the action from the map which topologically conjugates it to $`\mathrm{𝗉𝗋𝗈𝗃}`$.
The choice of the projective action as a “reference” compactification comes from the classification results (Theorems 1 and 2).
Let us prove Proposition 3.
### 1.1 First step : action on the boundary
We shall first prove that the restriction of $`\rho `$ to the boundary sphere is topologically conjugate to that of $`\mathrm{𝗉𝗋𝗈𝗃}`$.
Since the Hadamard boundary of the Poincaré ball is the only $`(n1)`$-spherical homogeneous space of $`G`$ (up to topological conjugacy), it is sufficient to prove that $`\rho `$ is transitive on the boundary.
Since $`\rho `$ is homeomorphic to $`\mathrm{𝗂𝗌𝗈𝗆}`$ in the open ball, the orbits in the open ball under the action $`\rho (\text{SO(n)})`$ (where SO(n) is seen as a subgroup of $`G`$) are: a fixed point $`O`$, and topological spheres disconnecting the open ball in two connected components. Moreover, $`O`$ is in the interior component of all these spherical orbits.
Take a point $`x`$ on the boundary of the closed ball and choose a real number $`\epsilon >0`$, smaller than the distance between the boundary and $`O`$.
By uniform continuity (SO(n) is compact), there is an orbit $`S`$ of $`\rho (\text{SO(n)})`$ and a point $`x^{}S`$ such that for all $`g\text{SO(n)}`$, $`\rho (g)x`$ is $`\epsilon `$-close to $`\rho (g)x^{}S`$. Therefore $`S`$ is in the neigborhood of size $`\epsilon `$ around the boundary. Since its interior contain $`O`$ (this condition is important, figure 2 shows a counterexample), any point of the boundary is $`\epsilon `$-close to $`S`$. Indeed, if $`p`$ is in the boundary and not $`\epsilon `$-close to $`S`$, there is a path avoiding $`S`$ and connecting $`p`$ to a point $`q`$ not $`\epsilon `$-close to the boundary. Since $`S`$ is in an $`\epsilon `$-neighborhood of the boundary, there is a path avoiding $`S`$ and connecting $`q`$ to $`O`$, thus there is a path avoiding $`S`$ and connecting $`p`$ (which is outside $`S`$) to $`O`$ (which is inside $`S`$), a contradiction.
The orbit of $`x`$ under $`\rho (\text{SO(n)})`$ is $`\epsilon `$-dense in the boundary for all sufficiently small $`\epsilon `$, thus it is dense. It is also compact, thus $`\rho `$ must be transitive on the boundary.
### 1.2 Second step : closure of a geodesic
Let $`\psi `$ be the homeomorphism of the open ball which conjugates $`\mathrm{𝗂𝗌𝗈𝗆}`$ (written in the Klein model) and the restriction of $`\rho `$ to the open ball.
Let $`L`$ be a projective geodesic. We shall prove that the closure of $`\psi (L)`$ in the closed ball has exactly two points in the boundary.
We denote by $`\text{EP(}L\text{)}`$ the intersection of the boundary and of the closure of $`\psi (L)`$.
Since $`\psi (L)`$ is non compact and closed in the open ball, $`\text{EP(}L\text{)}`$ is non empty. It is also compact and globally invariant under the action of $`\rho (g)`$ for all $`g`$ whose projective action leaves $`L`$ globally invariant. The projective action of these $`g`$’s has two orbits in the boundary : the first contains the two endpoints of $`L`$ and the second all the other points of the boundary. Thus, since the restriction of $`\rho `$ to the boundary is topologically conjugate to that of $`\mathrm{𝗉𝗋𝗈𝗃}`$, $`\text{EP(}L\text{)}`$ has two points or is the entire boundary (it is compact).
If $`\text{EP(}L\text{)}`$ were the entire boundary, let $`x`$ be a point in the boundary not fixed by some rotation $`\rho (g)`$ which leaves $`\psi (L)`$ pointwise invariant. Let $`(x_n)`$ be a sequence of points of $`\psi (L)`$ such that $`limx_n=x`$. Then $`\rho (g)x_n=x_n`$ thus by continuity $`\rho (g)x=x`$, a contradiction.
Therefore $`\psi (L)`$ has exactly two endpoints (one for each half of the geodesic). Moreover, the endpoints are determinated by their stabilizer, thus the images of two asymptotic half-geodesics have the same endpoint and any point of the boundary is the endpoint of some geodesic.
### 1.3 Third step : extension of $`\psi `$
We can extend $`\psi `$ into a one-to-one map (denoted by $`\stackrel{~}{\psi }`$) of the closed ball into itself : $`\stackrel{~}{\psi }`$ maps the endpoint of a half geodesic to the endpoint of its image by $`\psi `$. Then $`\stackrel{~}{\psi }`$ conjugates $`\rho `$ and $`\mathrm{𝗉𝗋𝗈𝗃}`$.
We shall prove that $`\stackrel{~}{\psi }`$ is a homeomorphism. Since it is one-to-one and the closed ball is compact, it is sufficient to prove its continuity. Since it coincides with $`\psi `$ in the interior, it is sufficient to prove its continuity at all points of the boundary.
Let $`x`$ be a point of the boundary and $`(x_n)_n`$ be a sequence of points of the closed ball such that $`limx_n=x`$.
Let $`L`$ be the closure of a half-geodesic of endpoint $`x`$. Let $`o`$ be a point of $`L`$, $`L_n`$ be the half-geodesic with endpoint $`o`$ and containing $`x_n`$, $`P_n`$ be the totally geodesic plane containing both $`L`$ and $`L_n`$ and $`Q_n`$ be the maximal totally geodesic subspace orthogonal to $`P_n`$. There is a sequence $`(g_n)_n`$ of elliptic elements of $`G`$ such that $`\mathrm{𝗉𝗋𝗈𝗃}_{g_n}`$ is identity on $`Q_n`$, $`\mathrm{𝗉𝗋𝗈𝗃}_{g_n}`$ globally stabilizes $`P_n`$ and $`\mathrm{𝗉𝗋𝗈𝗃}_{g_n}(L)=L_n`$ for all $`n`$. Let $`Y_n=\mathrm{𝗉𝗋𝗈𝗃}_{g_n^1}x_n`$. For all $`n`$, $`y_nL`$ and since $`x_n`$ has $`x`$ for limit, the sequence of angles of the rotations $`g_n`$ has limit zero and $`\underset{n\mathrm{}}{lim}g_n=1`$.
Then
$`\stackrel{~}{\psi }(x_n)`$ $`=`$ $`\stackrel{~}{\psi }(\mathrm{𝗉𝗋𝗈𝗃}_{g_n}y_n)`$
$`=`$ $`\rho (g_n)\stackrel{~}{\psi }(y_n)`$
and since $`y_n`$ has for limit the endpoint of $`L`$, by definition $`\stackrel{~}{\psi }(y_n)`$ has for limit the endpoint of $`\stackrel{~}{\psi }(L)`$, that is to say $`\stackrel{~}{\psi }(x)`$. Thus the uniform continuity of $`\rho `$ in some neighborhood of $`1G`$ ensures $`\underset{n\mathrm{}}{lim}\stackrel{~}{\psi }(x_n)=\stackrel{~}{\psi }(x)`$ and $`\stackrel{~}{\psi }`$ is continuous. The proof is complete.
## 2 Proof of Theorems 1 and 2
### 2.1 Main part
We shall start with a lemma which contains the heart of the proof. Recall that $`k`$ is always $`\mathrm{}`$ or $`\omega `$.
###### Lemma 5
Let $`\phi `$ be a $`𝒞^k`$ compactification of the isometric action of $`G`$ such that the inverse images of the projective geodesics are border-transversal $`𝒞^k`$ submanifolds of the closed ball. Then, up to a $`𝒞^k`$ change of coordinates, $`\phi `$ can be written in half-space charts (the chart at the goal being $`𝒦𝒞`$) in the following form:
$$\phi =(x_1,\mathrm{},x_{n1},y)(x_1,\mathrm{},x_{n1},f(y))$$
where $`f`$ is a $`𝒞^k`$ map.
We recall that for a map $`f`$ defined on $`^+`$, being $`𝒞^k`$ means that $`f`$ is $`𝒞^k`$ on $`^+^{}`$ and can be prolonged in a neighborhood of $`0`$ into a $`𝒞^k`$ map.
Proof: in the Lie algebra $`𝔤`$ of $`G`$ we shall denote by $`X_1,X_2,\mathrm{},X_{n1}`$ a basis of the vector space of the parabolic transformations which fix $`\mathrm{}`$ (the only point missed by the half-space chart). The corresponding vector fields for the projective action are denoted by $`\mathrm{𝗉𝗋𝗈𝗃}_{X_1},\mathrm{},\mathrm{𝗉𝗋𝗈𝗃}_{X_{n1}}`$. By definition, in the chart $`𝒦𝒞`$ they may be written as $`\mathrm{𝗉𝗋𝗈𝗃}_{X_i}=\frac{}{x_i}`$.
Without loss of generality, we can suppose that $`\phi (0)=0`$, $`\phi (\mathrm{})=\mathrm{}`$ and, choosing wisely the $`\overline{}^{n+}`$ chart for the source of $`\phi `$, that $`\phi {}_{}{}^{}(\mathrm{𝗉𝗋𝗈𝗃}_{X_i})=\frac{}{x_i}=\mathrm{𝗉𝗋𝗈𝗃}_{X_i}`$ for each $`i`$ (the subgroup generated by $`\phi {}_{}{}^{}(\mathrm{𝗉𝗋𝗈𝗃}_{X_i})`$’s acts freely and is abelian.)
Note that in the chart $`𝒦𝒞`$ any vertical line is a geodesic.
We shall first prove that via a differentiable change of coordinates we can suppose that the inverse image of the $`y`$ axis is the $`y`$ axis itself. Let $`L`$ be the inverse image of the $`y`$ axis (that is to say the geodesic joining 0 to $`\mathrm{}`$). Then $`L`$ meets each horizontal hyperplane only once (these hyperplanes are the orbits of the action of the $`\phi {}_{}{}^{}(\mathrm{𝗉𝗋𝗈𝗃}_{X_i})`$’s) and by hypothesis is a border-transversal $`𝒞^k`$ submanifold of the closed half space. Hence
$$L=\left\{(f_1(y),\mathrm{},f_n(y),y);y^+\right\}$$
where $`f_i`$’s are $`𝒞^k`$ maps. The differentiable change of coordinates
$$(x_1,\mathrm{},x_n,y)(x_1f_1(y),\mathrm{},x_nf_n(y),y)$$
transforms $`L`$ into the $`y`$ axis and do not change $`\phi {}_{}{}^{}(\mathrm{𝗉𝗋𝗈𝗃}_{X_i})`$.
When the inverse image of the $`y`$ axis is the $`y`$ axis itself, there is a continuous map $`f:^+^+`$ such that
$$\phi (0,\mathrm{},0,y)=(0,\mathrm{},0,f(y)).$$
Since $`\phi {}_{}{}^{}(\frac{}{x_i})=\frac{}{x_i}`$ we have
$$\phi (x_1,\mathrm{},x_n,y)=(x_1,\mathrm{},x_n,f(y)).$$
There is a projective geodesic which may be written at 0 as
$$\left\{(x,k_2(x),\mathrm{},k_n(x),k_{n+1}(x));x[0,\epsilon ]\right\}$$
where $`k_i`$’s are $`𝒞^k`$ maps; let $`L_2`$ be its inverse image.
By hypothesis, we can write $`L_2=\left\{(l_1(y),\mathrm{},l_n(y),y);y^+\right\}`$ where $`l_i`$’s are $`𝒞^k`$ maps. Then computing $`\phi (L_2)=\left\{(l_1(y),\mathrm{},l_n(y),f(y));y^+\right\}`$ gives $`f=k_{n+1}l_1`$, hence $`f`$ is $`𝒞^k`$.
### 2.2 Transversality of geodesics in dimension at least 3
###### Lemma 6
Let $`\phi `$ be a $`𝒞^k`$ compactification of $`\mathrm{𝗂𝗌𝗈𝗆}`$. If $`n3`$ then the inverse images of the projective geodesics are $`𝒞^k`$ border-transversal submanifolds.
Proof: If $`n`$ is odd, each geodesic is the set of fixed points of a symmetry whose differentials in the endpoints of the geodesic has the form $`\left(\begin{array}{cccc}1& & & \\ & 1& & \\ & & \mathrm{}& \\ & & & 1\end{array}\right)`$.
If $`n`$ is even, each geodesic is the set of common fixed points of two symmetries around totally geodesic planes whose differentials in the endpoints of the geodesic have the form $`\left(\begin{array}{ccccc}1& & & & \\ & 1& & & \\ & & 1& & \\ & & & \mathrm{}& \\ & & & & 1\end{array}\right)`$ and $`\left(\begin{array}{ccccc}1& & & & \\ & 1& & & \\ & & \mathrm{}& & \\ & & & 1& \\ & & & & 1\end{array}\right)`$.
In both cases, we get implicit $`𝒞^k`$ definitions of the geodesics, which are therefore border-transversal $`𝒞^k`$ submanifolds of the closed ball.
With lemmas 5 and 6 we have proved that any compactification is given by a map $`\phi _f`$ of the form (2) for some $`𝒞^k`$ map $`f`$.
It is clear that two such compactifications are conjugate if and only of the corresponding maps $`f`$’s are.
### 2.3 End of proof of Theorem 1
Concerning Theorem 1 we are left with proving that we can replace the map $`f`$ by a monomial map and that any monomial map gives a compactification.
We omit the proof of the following classical fact.
###### Lemma 7
In the analytic case, we can replace $`f`$ by a map of the form $`yy^p`$
###### Lemma 8
Let $`\overline{X}`$ be an analytic vector field on the closed half-plane, tangent to the boundary. Then for any integer $`p`$, the pull-back of $`\overline{X}`$ by $`(x_1,\mathrm{},x_{n1},y)(x_1,\mathrm{},x_{n1},y^p)`$ extends analytically to the boundary.
Proof: since $`\overline{X}`$ is analytic, it has the form
$$\overline{X}=\underset{i=1}{\overset{n}{}}\underset{a,b}{}\alpha _{a,b}^ix^ay^b\frac{}{x_i}+\underset{a,b}{}\beta _{a,b}x^ay^b\frac{}{y}$$
where the sums are taken over all non-negative integers $`b`$ and all $`n`$-tuples of non-negative integers $`a`$; $`x^a`$ means $`x_1^{a_1}x_2^{a_2}\mathrm{}x_n^{a_n}`$.
We denote by $`\phi _p`$ the map $`(x_1,\mathrm{},x_{n1},y)(x_1,\mathrm{},x_{n1},y^p)`$. After a direct computation, we see that
$$\phi _p{}_{}{}^{}(\overline{X})=\underset{i=1}{\overset{n}{}}\underset{a,b}{}\alpha _{a,b}^ix^ay^{pb}\frac{}{x_i}+\underset{a,b}{}\beta _{a,b}x^ay^{pb+1p}\frac{}{y}$$
and hence, it is analytic if for all $`a`$, $`\beta _{a,0}=0`$, that is to say if $`\overline{X}`$ is tangent to the boundary.
###### Corollary 9
Let $`p`$ be a positive integer. Then the map $`\phi _p`$ given by (1) is a compactification.
Proof: Lemma 8 ensures that the pull-back of the action of the Lie algebra of $`G`$ admits an extension to the closed ball into an analytic action. This action is complete by compacity, thus it gives an action of the universal cover $`\stackrel{~}{G}`$ of $`G`$. Let $`g`$ be an element $`\stackrel{~}{G}`$ which projects on identity. Then $`g`$ acts trivially in the open ball, therefore it acts trivially in the whole closed ball. Thus, the action of $`𝔤`$ gives an action of $`G`$.
### 2.4 End of proof of Theorem 2
Concerning Theorem 2 we are left with proving that the map $`\phi _f`$ given by (2) is a compactification if and only if $`f`$ satisfies (3) and that this condition is satisfied by all non-flat maps.
###### Lemma 10
In any dimension $`n2`$, $`\phi _f`$ is a smooth compactification if and only if $`f`$ satisfies (3)
Proof: As seen in the proof of 9, $`\phi `$ gives a compactification of $`\mathrm{𝗂𝗌𝗈𝗆}`$ if and only if the vector fields $`\phi {}_{}{}^{}(\mathrm{𝗉𝗋𝗈𝗃}_X)`$ admits a $`𝒞^{\mathrm{}}`$ extension to the boundary.
Let $`H`$ be a hyperbolic element of $`𝔤`$ stabilizing the $`y`$ axis, $`(X_i)_{1in1}`$ be a basis of the vector space of the parabolic elements of $`𝔤`$ stabilizing $`\mathrm{}`$ (the only point not contained in the chart $`𝒦𝒞`$), $`(Y_i)_{1in1}`$ be a basis of the vector space of the parabolic elements of $`𝔤`$ stabilizing $`0`$, and $`(R_j)_{1j\frac{(n1)(n2)}{2}}`$ be a basis of the vector space of the elliptic elements of $`𝔤`$ stabilizing (pointwise) the $`y`$ axis.
The union of this elements generates $`𝔤`$, thus we only have to check that they all admit smooth extensions in the model given by $`\phi `$.
The $`\mathrm{𝗉𝗋𝗈𝗃}_{X_i}`$’s and $`\mathrm{𝗉𝗋𝗈𝗃}_{R_j}`$’s are left unchanged by $`\phi `$, thus admit extensions.
The $`Y_i`$’s are conjugate one to another by rotations of $`G`$ generated by the $`R_j`$’s, so we just have to check one of them.
We can explicitly compute $`\stackrel{~}{H}=\phi {}_{}{}^{}(\mathrm{𝗉𝗋𝗈𝗃}_H)`$ and one of the $`\stackrel{~}{Y}_i=\phi {}_{}{}^{}(\mathrm{𝗉𝗋𝗈𝗃}_{Y_i})`$. We first give the explicit expressions of $`\mathrm{𝗉𝗋𝗈𝗃}_H`$ and $`\mathrm{𝗉𝗋𝗈𝗃}_{Y_1}`$ in the chart $`𝒦𝒞`$, computed thanks to (4).
$`\mathrm{𝗉𝗋𝗈𝗃}_H(x_1,\mathrm{},x_{n1},y)`$ $`=`$ $`|\begin{array}{c}2x_1\hfill \\ 2x_2\hfill \\ \mathrm{}\hfill \\ 2x_{n1}\hfill \\ 4y\hfill \end{array}`$ (10)
$`\mathrm{𝗉𝗋𝗈𝗃}_{Y_1}(x_1,\mathrm{},x_{n1},y)`$ $`=`$ $`|\begin{array}{c}y+x_2^2+x_3^2+\mathrm{}+x_{n1}^2x_1^2\hfill \\ 2x_1x_2\hfill \\ 2x_1x_3\hfill \\ \mathrm{}\hfill \\ 2x_1x_{n1}\hfill \\ 4x_1y\hfill \end{array}`$ (17)
$`\stackrel{~}{H}(x_1,\mathrm{},x_{n1},y)`$ $`=`$ $`|\begin{array}{c}2x_1\hfill \\ 2x_2\hfill \\ \mathrm{}\hfill \\ 2x_{n1}\hfill \\ 4f(y)/f^{}(y)\hfill \end{array}`$ (23)
$`\stackrel{~}{Y}_1(x_1,\mathrm{},x_{n1},y)`$ $`=`$ $`|\begin{array}{c}f(y)+x_2^2+x_3^2+\mathrm{}+x_{n1}^2x_1^2\hfill \\ 2x_1x_2\hfill \\ 2x_1x_3\hfill \\ \mathrm{}\hfill \\ 2x_1x_{n1}\hfill \\ 4x_1f(y)/f^{}(y)\hfill \end{array}`$ (30)
We see that $`\stackrel{~}{H}`$ and $`\stackrel{~}{Y}_1`$ admit smooth extensions if and only if $`f/f^{}`$ does.
###### Lemma 11
Let $`f`$ be a smooth homeomorphism of $`^+`$. If $`f`$ is non-flat, then it satisfies (3).
We omit the proof, a simple verification.
UMPA, ÉNS Lyon
46, allée d’Italie
69 364 Lyon cedex 07
France
www.umpa.ens-lyon.fr/$``$bkloeckn/
bkloeckn@umpa.ens-lyon.fr |
warning/0506/gr-qc0506113.html | ar5iv | text | # 1 Introduction
## 1 Introduction
The standard Friedmann-Robertson-Walker (FRW) cosmological model prescribes a homogeneous and an isotropic distribution for its matter in the description of the present state of the universe. At the present state of evolution, the universe is spherically symmetric and the matter distribution in the universe is on the whole isotropic and homogeneous. But in early stages of evolution, it could have not had such a smoothed picture. Close to the big bang singularity, neither the assumption of spherical symmetry nor that of isotropy can be strictly valid. So we consider plane-symmetric, which is less restrictive than spherical symmetry and can provide an avenue to study inhomogeneities. Inhomogeneous cosmological models play an important role in understanding some essential features of the universe such as the formation of galaxies during the early stages of evolution and process of homogenization. The early attempts at the construction of such models have done by Tolman and Bondi who considered spherically symmetric models. Inhomogeneous plane-symmetric models were considered by Taub and later by Tomimura, Szekeres, Collins and Szafron, Szafron and Collins. Recently, Senovilla obtained a new class of exact solutions of Einstein’s equation without big bang singularity, representing a cylindrically symmetric, inhomogeneous cosmological model filled with perfect fluid which is smooth and regular everywhere satisfying energy and causality conditions. Later, Ruis and Senovilla have separated out a fairly large class of singularity free models through a comprehensive study of general cylindrically symmetric metric with separable function of $`r`$ and $`t`$ as metric coefficients. Dadhich et al. have established a link between the FRW model and the singularity free family by deducing the latter through a natural and simple in-homogenization and anisotropization of the former. Recently, Patel et al. presented a general class of inhomogeneous cosmological models filled with non-thermalized perfect fluid by assuming that the background space-time admits two space-like commuting killing vectors and has separable metric coefficients. Bali and Tyagi obtained a plane-symmetric inhomogeneous cosmological models of perfect fluid distribution with electro-magnetic field. Recently, Pradhan et al. have investigated a plane-symmetric inhomogeneous viscous fluid cosmological models with electro-magnetic field.
Models with a relic cosmological constant $`\mathrm{\Lambda }`$ have received considerable attention recently among researchers for various reasons (see Refs.$``$ and references therein). Some of the recent discussions on the cosmological constant “problem” and consequence on cosmology with a time-varying cosmological constant by Ratra and Peebles, Dolgov$``$ and Sahni and Starobinsky have pointed out that in the absence of any interaction with matter or radiation, the cosmological constant remains a “constant”. However, in the presence of interactions with matter or radiation, a solution of Einstein equations and the assumed equation of covariant conservation of stress-energy with a time-varying $`\mathrm{\Lambda }`$ can be found. For these solutions, conservation of energy requires decrease in the energy density of the vacuum component to be compensated by a corresponding increase in the energy density of matter or radiation. Earlier researchers on this topic, are contained in Zeldovich Weinberg and Carroll, Press and Turner. Recent observations by Perlmutter et al. and Riess et al. strongly favour a significant and positive value of $`\mathrm{\Lambda }`$. Their finding arise from the study of more than $`50`$ type Ia supernovae with redshifts in the range $`0.10z0.83`$ and these suggest Friedman models with negative pressure matter such as a cosmological constant $`(\mathrm{\Lambda })`$, domain walls or cosmic strings (Vilenkin , Garnavich et al.). Recently, Carmeli and Kuzmenko have shown that the cosmological relativistic theory (Behar and Carmeli) predicts the value for cosmological constant $`\mathrm{\Lambda }=1.934\times 10^{35}s^2`$. This value of “$`\mathrm{\Lambda }`$” is in excellent agreement with the measurements recently obtained by the High-Z Supernova Team and Supernova Cosmological Project (Garnavich et al. ; Perlmutter et al.; Riess et al.; Schmidt et al.). The main conclusion of these observations is that the expansion of the universe is accelerating.
Several ans$`\ddot{a}`$tz have been proposed in which the $`\mathrm{\Lambda }`$ term decays with time (see Refs. Gasperini, Berman, Freese et al., $`\ddot{O}`$zer and Taha, Peebles and Ratra, Chen and Hu, Abdussattar and Viswakarma, Gariel and Le Denmat, Pradhan et al.). Of the special interest is the ans$`\ddot{a}`$tz $`\mathrm{\Lambda }S^2`$ (where $`S`$ is the scale factor of the Robertson-Walker metric) by Chen and Wu, which has been considered/modified by several authors ( Abdel-Rahaman, Carvalho et al., Waga, Silveira and Waga, Vishwakarma).
Most cosmological models assume that the matter in the universe can be described by ’dust’ (a pressure-less distribution) or at best a perfect fluid. However, bulk viscosity is expected to play an important role at certain stages of expanding universe$``$. It has been shown that bulk viscosity leads to inflationary like solution, and acts like a negative energy field in an expanding universe. Furthermore, there are several processes which are expected to give rise to viscous effects. These are the decoupling of neutrinos during the radiation era and the decoupling of radiation and matter during the recombination era. Bulk viscosity is associated with the Grand Unification Theories (GUT) phase transition and string creation. Thus, we should consider the presence of a material distribution other than a perfect fluid to have realistic cosmological models (see Grøn for a review on cosmological models with bulk viscosity). A number of authors have discussed cosmological solutions with bulk viscosity in various context$``$.
The occurrence of magnetic fields on galactic scale is well-established fact today, and their importance for a variety of astrophysical phenomena is generally acknowledged as pointed out Zeldovich et al.. Also Harrison has suggested that magnetic field could have a cosmological origin. As a natural consequences, we should include magnetic fields in the energy-momentum tensor of the early universe. The choice of anisotropic cosmological models in Einstein system of field equations leads to the cosmological models more general than Robertson-Walker model. The presence of primordial magnetic fields in the early stages of the evolution of the universe has been discussed by several authors$``$. Strong magnetic fields can be created due to adiabatic compression in clusters of galaxies. Large-scale magnetic fields give rise to anisotropies in the universe. The anisotropic pressure created by the magnetic fields dominates the evolution of the shear anisotropy and it decays slower than if the pressure was isotropic. Such fields can be generated at the end of an inflationary epoch$``$. Anisotropic magnetic field models have significant contribution in the evolution of galaxies and stellar objects. The model studied by Murphy possessed an interesting feature is that the big bang type of singularity of infinite spacetime curvature does not occur to a finite past. However, the relationship assumed by Murphy between the viscosity coefficient and the matter density is not acceptable at large density. Several authors $``$ have investigated cosmological models with a magnetic field in different context.
Recently Bali et al. obtained some plane-symmetric inhomogeneous cosmological models for a perfect fluid distribution with electro-magnetic field. Motivated the situations discussed above, in this paper, we shall focus upon the problem of establishing a formalism for studying the general relativistic evolution magnetic inhomogeneities in presence of bulk viscous in an expanding universe. We do this by extending the work of Bali et al. by including an electrically neutral bulk viscous fluid as the source of matter in the energy-momentum tensor. This paper is organized as follows. The metric and the field equations are presented in section 2. In section 3 we deal with the solution of the field equations in presence of bulk viscous fluid. The sections 3.1, 3.2 and 3.3 contain the two different cases (i.e. for $`n=0`$ and $`n=1`$ ) and also contain some physical aspects of these models respectively. Section $`4`$ describe some other generated models and their physical and geometric properties. Finally in section $`5`$ concluding remarks have been given.
## 2 The metric and field equations
We consider the metric in the form
$$ds^2=B^2(dx^2dt^2+dy^2)+C^2dz^2,$$
(1)
where the metric potential $`B`$ and $`C`$ are functions of $`x`$ and $`t`$.
The energy momentum tensor in the presence of bulk stress has the form
$$T_i^j=(\rho +\overline{p})v_iv^j+\overline{p}g_i^j+E_i^j,$$
(2)
where $`E_i^j`$ is the electro-magnetic field given by
$$E_i^j=F_{i\alpha }F^{j\alpha }\frac{1}{4}F_{\alpha \beta }F^{\alpha \beta }g_i^j,$$
(3)
and
$$\overline{p}=p\xi v_{;i}^i.$$
(4)
Here $`\rho `$, $`p`$, $`\overline{p}`$ , $`F_i^j`$ and $`\xi `$ are the energy density, isotropic pressure, effective pressure, electromagnetic field tensor and bulk viscous coefficient respectively and $`v^i`$ is the flow vector satisfying the relation
$$g_{ij}v^iv^j=1.$$
(5)
Here the semicolon represents a covariant differentiation. The coordinates are considered to be comoving so that $`v^1`$ = $`0`$ = $`v^2`$ = $`v^3`$ and $`v^4`$ = $`\frac{1}{B}`$. We consider the current to be flowing along the $`z`$-axis so that $`F_{12}`$ is the only non-vanishing component of $`F_{ij}`$.
The Einstein’s field equations ( in gravitational units c = 1, G = 1 ) read as
$$R_i^j\frac{1}{2}Rg_i^j+\mathrm{\Lambda }g_i^j=8\pi T_i^j,$$
(6)
for the line element (1) has been set up as
$$8\pi B^2\left(\overline{p}+\frac{F_{12}^2}{2B^4}\right)=\frac{B_{44}}{B}+\frac{B_4^2}{B^2}+\frac{B_1^2}{B^2}\frac{C_{44}}{C}+\frac{2B_1C_1}{BC}\mathrm{\Lambda }B^2,$$
(7)
$$8\pi B^2\left(\overline{p}+\frac{F_{12}^2}{2B^4}\right)=\frac{B_{44}}{B}+\frac{B_4^2}{B^2}+\frac{B_{11}}{B}\frac{B_1^2}{B^2}\frac{C_{44}}{C}+\frac{C_{11}}{C}\mathrm{\Lambda }B^2,$$
(8)
$$8\pi B^2\left(\overline{p}\frac{F_{12}^2}{2B^4}\right)=\frac{2B_{44}}{B}+\frac{B_4^2}{B^2}+\frac{2B_{11}}{B}\frac{B_1^2}{B^2}\mathrm{\Lambda }B^2,$$
(9)
$$8\pi B^2\left(\rho +\frac{F_{12}^2}{2B^4}\right)=\frac{B_{11}}{B}\frac{C_{11}}{C}+\frac{B_1^2}{B^2}+\frac{2B_4^2}{B^2}+\frac{2B_4C_4}{BC}+\mathrm{\Lambda }B^2,$$
(10)
$$0=\frac{B_{14}}{B}+\frac{C_{14}}{C}\frac{2B_1B_4}{B^2}\frac{B_1C_4}{BC}\frac{B_4C_1}{BC},$$
(11)
where
$$\overline{p}=p\frac{\xi }{B}\left(\frac{2B_4}{B}+\frac{C_4}{C}\right).$$
Here and in the following expressions the suffixes $`1`$ and $`4`$ at the symbols $`B`$, $`C`$, $`f`$ and $`g`$ denote differentiation with respect to $`x`$ and $`t`$ respectively.
## 3 Solution of the field equations
From Equations (7), (8) and (9), we have
$$\frac{B_{11}}{B}+\frac{C_{11}}{C}\frac{2B_1C_1}{BC}\frac{2B_1^2}{B^2}=0,$$
(12)
$$\frac{8\pi F_{12}^2}{B^2}=\frac{B_{44}}{B}\frac{B_{11}}{B}+\frac{C_{11}}{C}\frac{C_{44}}{C}.$$
(13)
Equations (7)-(11) represent a system of five equations in seven unknowns $`B`$, $`C`$, $`\rho `$, $`p`$, $`F_{12}`$, $`\mathrm{\Lambda }`$ and $`\xi `$. The research on exact solutions is based on some physically reasonable restrictions used to simplify the Einstein equations. To get a determinate solution, we need two extra conditions.
Let us consider that
$$B=f(x)g(t),$$
$$C=h(x)k(t).$$
(14)
Using Equation (14) in Equation (11), we obtain
$$\frac{\frac{f_1}{f}}{\frac{h_1}{h}}=\frac{\frac{k_4}{k}\frac{g_4}{g}}{\frac{k_4}{k}+\frac{g_4}{g}}=K(constant,say),$$
(15)
which leads to
$$\frac{f_1}{f}=K\frac{h_1}{h},$$
(16)
and
$$\frac{k_4}{k}\frac{g_4}{g}=K\left(\frac{k_4}{k}+\frac{g_4}{g}\right).$$
(17)
Using Equation (14) in Equation (12) leads to
$$\frac{f_{11}}{f}\frac{2f_1}{f}\left(\frac{h_1}{h}+\frac{f_1}{f}\right)+\frac{h_{11}}{h}=0.$$
(18)
Equations (16) and (18) give
$$(K+1)hh_{11}(3K+K^2)h_1^2=0,$$
(19)
which on integration leads to
$$h=\left[\frac{K+1}{(12KK^2)(\alpha x+\beta )}\right]^{\frac{(K+1)}{(K^2+2K1)}},$$
(20)
where $`\alpha `$ and $`\beta `$ are constants of integration. Integrating Equations (16) and (17), we obtain
$$f=ah^K,$$
(21)
and
$$k=\left(\frac{g}{b}\right)^{\frac{(1+K)}{(1K)}}.$$
(22)
respectively, where $`a`$ and $`b`$ are constants of integration. Hence from Equations (14), (20)-(22), we obtain
$$B=ag\left[\frac{K+1}{(12KK^2)(\alpha x+\beta )}\right]^{\frac{K(K+1)}{(K^2+2K1)}},$$
(23)
and
$$C=\left(\frac{g}{b}\right)^{\frac{(1+K)}{(1K)}}\left[\frac{K+1}{(12KK^2)(\alpha x+\beta )}\right]^{\frac{(K+1)}{(K^2+2K1)}}.$$
(24)
After suitable transformation of coordinates and by taking $`\alpha `$ as unity without any loss of generality the metric (1) reduces to the form
$$ds^2=M^2\left(\frac{1}{X^2}\right)^{\frac{K(K+1)}{(K^2+2K1)}}g^2(T)\left(dX^2dT^2+dY^2\right)+$$
$$\left(\frac{1}{X^2}\right)^{\frac{(K+1)}{(K^2+2K1)}}\left\{g^2(T)\right\}^{\frac{1+K}{1K}}dZ^2,$$
(25)
where
$$M=a\left[\frac{K+1}{(12KK^2)}\right]^{\frac{K(K+1)}{(K^2+2K1)}}.$$
(26)
There is a lot of known solutions to the Einstein field equations but Equation (25) is indeed a new one. The effective pressure $`\overline{p}`$ and density $`\rho `$ for the model (25) are given by
$$8\pi \overline{p}=\mathrm{\Lambda }+\frac{1}{g^2M^2X^{\frac{2K(K+1)}{(12KK^2)}}}\times $$
$$\left[\frac{2K^2(K+1)(K+2)}{(12KK^2)^2X^2}\frac{(2K)}{(1K)}\frac{g_{44}}{g}+\frac{(13K)}{(1K)^2}\frac{g_4^2}{g^2}\right],$$
(27)
$$8\pi \rho =\mathrm{\Lambda }+\frac{1}{g^2M^2X^{\frac{2K(K+1)}{(12KK^2)}}}\times $$
$$\left[\frac{2K(K+1)(K+2)}{(12KK^2)^2X^2}\frac{1}{(1K)}\frac{g_{44}}{g}+\frac{(45KK^2)}{(1K)^2}\frac{g_4^2}{g^2}\right].$$
(28)
For the specification of $`\xi `$, we assume that the fluid obeys an equation of state of the form
$$p=\gamma \rho ,$$
(29)
where $`\gamma (0\gamma 1)`$ is a constant.
Thus, given $`\xi (t)`$ we can solve for the cosmological parameters. In most of the investigations involving bulk viscosity is assumed to be a simple power function of the energy density$``$.
$$\xi (t)=\xi _0\rho ^n,$$
(30)
where $`\xi _0`$ and $`n`$ are constants. If $`n=1`$, Equation (26) may correspond to a radiative fluid. However, more realistic models are based on $`n`$ lying in the regime $`0n\frac{1}{2}`$.
On using (30) in (27), we obtain
$$8\pi (p\xi _0\rho ^n\theta )=\mathrm{\Lambda }+\frac{1}{g^2M^2X^{\frac{2K(K+1)}{(12KK^2)}}}\times $$
$$\left[\frac{2K^2(K+1)(K+2)}{(12KK^2)^2X^2}\frac{(2K)}{(1K)}\frac{g_{44}}{g}+\frac{(13K)}{(1K)^2}\frac{g_4^2}{g^2}\right],$$
(31)
where $`\theta `$ is the scalar of expansion calculated for the flow vector $`v^i`$ and is given by
$$\theta =\frac{g_4}{g^2\sqrt{M}X^{\frac{K(K+1)}{(12KK^2)}}}\left(\frac{K3}{K1}\right).$$
(32)
### 3.1 Model I: $`(\xi =\xi _0)`$
When $`n=0`$, Equation (30) reduces to $`\xi =\xi _0`$. With the use of Equations (28), (29) and (32), Equation (31) reduces to
$$8\pi (1+\gamma )\rho =\frac{8\pi \xi _0g_4}{g^2\sqrt{M}X^{\frac{K(K+1)}{(12KK^2)}}}\left(\frac{K3}{K1}\right)+$$
$$\frac{1}{g^2M^2X^{\frac{2K(K+1)}{(12KK^2)}}}\left[\frac{2K(K^21)(K+2)}{(12KK^2)^2X^2}\frac{(3K)}{(1K)}\frac{g_{44}}{g}+\frac{(58KK^2)}{(1K)^2}\frac{g_4^2}{g^2}\right].$$
(33)
Eliminating $`\rho (t)`$ between (28) and (33), we get
$$(1+\gamma )\mathrm{\Lambda }=\frac{8\pi \xi _0g_4}{g^2\sqrt{M}X^{\frac{K(K+1)}{(12KK^2)}}}\left(\frac{K3}{K1}\right)+\frac{1}{g^2M^2X^{\frac{2K(K+1)}{(12KK^2)}}}\times $$
$$\left[\frac{2K(K+1)(K+2)(K+\gamma )}{(12KK^2)^2X^2}\frac{(2K\gamma )}{(1K)}\frac{g_{44}}{g}+\frac{\left\{13K(45K+K^2)\gamma \right\}}{(1K^2)}\frac{g_4^2}{g^2}\right].$$
(34)
### 3.2 Model II: $`(\xi =\xi _0\rho )`$
When $`n=1`$, Equation (30) reduces to $`\xi =\xi _0\rho `$. With the use of (28), (29) and (32), Equation (31) reduces to
$$8\pi \left[1+\gamma \frac{\xi _0g_4}{g^2\sqrt{M}X^{\frac{K(K+1)}{(12KK^2)}}}\left(\frac{K3}{K1}\right)\right]\rho =$$
$$\frac{1}{g^2M^2X^{\frac{2K(K+1)}{(12KK^2)}}}\left[\frac{2K(K^21)(K+2)}{(12KK^2)^2X^2}\frac{(3K)}{(1K)}\frac{g_{44}}{g}+\frac{(58KK^2)}{(1K)^2}\frac{g_4^2}{g^2}\right].$$
(35)
Eliminating $`\rho (t)`$ between (28) and (35), we get
$$[1+\gamma \frac{\xi _0g_4}{g^2\sqrt{M}X^{\frac{K(K+1)}{(12KK^2)}}}\left(\frac{K3}{K1}\right)]\mathrm{\Lambda }=\frac{1}{g^2M^2X^{\frac{2K(K+1)}{(12KK^2)}}}\times $$
$$\left[\frac{2K(K+1)(K+2)(K+\gamma )}{(12KK^2)^2X^2}\frac{1\gamma (2K)}{(1K)}\frac{g_{44}}{g}+\frac{(45KK^2)(13K)\gamma }{(1K)^2}\frac{g_4^2}{g^2}\right]$$
$$+\frac{\xi _0g_4}{g^4M^{\frac{3}{2}}X^{\frac{3K(K+1)}{(12KK^2)}}}\left(\frac{K3}{K1}\right)\times $$
$$\left[\frac{2K(K+1)(K+2)}{(12KK^2)^2X^2}\frac{(2K)}{(1K)}\frac{g_{44}}{g}+\frac{(13K)}{(1K)^2}\frac{g_4^2}{g^2}\right].$$
(36)
In spite of homogeneity at large scale our universe is inhomogeneous at small scales, so physical quantities having position dependent are more natural in our observable universe if we do not go to super high scale. This result shows this kind of physical importance. In recent times the $`\mathrm{\Lambda }`$-term has interested theoreticians and observers for various reasons. The nontrivial role of the vacuum in the early universe generate a $`\mathrm{\Lambda }`$-term that leads to inflationary phase. Observationally this term provides an additional parameter to accommodate conflicting data on the values of the Hubble constant, the deceleration parameter, the density parameter and the age of the universe (for example, see the references ,). Assuming that $`\mathrm{\Lambda }`$ owes its origin to vacuum interactions, as suggested in particular by Sakharov it follows that it would in general be a function of space and time coordinates, rather than a strict constant. In a homogeneous universe $`\mathrm{\Lambda }`$ will be at most time dependent . In our case this approach can generate $`\mathrm{\Lambda }`$ that varies both with space and time. In considering the nature of local massive objects, however, the space dependence of $`\mathrm{\Lambda }`$ cannot be ignored. For details discussion, the readers are advised to see the references (Narlikar, Pecker and Vigier , Ray et al. ).
The effect of bulk viscosity is to introduce a change in the perfect fluid model. We also observe here that the condition of Murphy about the absence of a big bang type of singularity in the finite past in models with bulk viscous fluid is, in general, not true. We have freedom of choosing the function $`g(T)`$ so that to give a physical behaviour of above parameters. As a matter of fact, there are multiple choices, for example, $`g(T)=c^2+d^2t^2,c^2+e^{dt^2},c^2+d^2\mathrm{cos}\omega t,c^2>d^2`$, where $`c`$ and $`d`$ are some real constants. From Equations (34) and (36), we observe that the cosmological constant is a decreasing function of time and it approaches a small positive value at late times under some suitable conditions which explains the small value of $`\mathrm{\Lambda }`$ at present.
### 3.3 Some physical aspects of the models
We shall now give the expressions for kinematical quantities and the components of conformal curvature tensor. With regard to the kinematical properties of the velocity vector $`v^i`$ in the metric (25), a straightforward calculation leads to the following expression for the shear of the fluid:
$$\sigma ^2=\frac{4K^2g_{4}^{}{}_{}{}^{2}}{3\sqrt{M}g^4(1K)^2X^{\frac{2K(K+1)}{(12KK^2)}}}.$$
(37)
The rotation $`\omega `$ and acceleration are identically zero. The expansion scalar $`\theta `$ has already been given by (32). The non-vanishing physical components of the conformal curvature tensor are given by
$$C_{(1212)}=C_{(3434)}=\frac{1}{6Mg^2X^{\frac{2K(K+1)}{(12KK^2)}}}\times $$
$$\left[\frac{2K(K+1)(K1)}{(12KK^2)^2X^2}\frac{4K}{(1K)}\frac{g_{44}}{g}+\frac{4K(13K)}{(1K)^2}\frac{g_4^2}{g^2}\right],$$
(38)
$$C_{(1313)}=C_{(2424)}=\frac{1}{6Mg^2X^{\frac{2K(K+1)}{(12KK^2)}}}\times $$
$$\left[\frac{4K(K+1)(1K)}{(12KK^2)^2X^2}+\frac{2K}{(1K)}\frac{g_{44}}{g}+\frac{2K(3K1)}{(1K)^2}\frac{g_4^2}{g^2}\right],$$
(39)
$$C_{(1224)}=C_{(1334)}=\frac{K(K+1)}{M(12KK^2)X^{\frac{2K(K+1)}{(12KK^2)}}}\frac{g_4}{g},$$
(40)
$$C_{(2323)}=C_{(1414)}=\frac{1}{3Mg^2X^{\frac{2K(K+1)}{(12KK^2)}}}\left[\frac{K(K+1)^2(1K)}{(12KK^2)^2X^2}\frac{2K}{(1K)}\frac{g_4^2}{g^2}\right].$$
(41)
The non-vanishing component $`F_{12}`$ of the electromagnetic field tensor and $`J^2`$, the component of charge current density, are given by
$$F_{12}=\frac{1}{2\sqrt{\pi }}\left[\frac{K}{(1K)}\frac{g_{44}}{g}\frac{K(1+K)}{(1K)^2}\frac{g_4^2}{g^2}\frac{K(K+1)(K1)(K+2)}{(12KK^2)X^2}\right]^{\frac{1}{2}}B,$$
(42)
$$J^2=\frac{1}{M^2g^4X^{\frac{(1+K)(1+3K)}{(12KK^2)}}}\times $$
$$\left[\frac{\sqrt{\zeta }(1K^2)}{(12KK^2)}X^{\frac{2K}{(12KK^2)}}+\frac{2K(K+1)(K1)(K+2)\zeta ^{\frac{1}{2}}}{(12KK^2)^2}X^{\frac{2(K^2+3K1)}{(12KK^2)}}\right],$$
(43)
where
$$\zeta =\left[\frac{2K}{(1K)}\frac{g_{44}}{g}\frac{2K(1+K)}{(1K)^2}\frac{g_4^2}{g^2}\frac{2K(K+1)(K1)(K+2)}{(12KK^2)^2X^2}\right].$$
(44)
The models represent shearing, non-rotating and Petrov type I non-degenerate in general, in which the flow is geodetic. The expansion in the model stops when $`K=3`$ but it will continue indefinitely when $`K>3`$. Since $`lim_{T0}\frac{\sigma }{\theta }0`$, hence the models do not approach isotropy for large values of $`T`$.
## 4 Other generated model
In the metric (25), the function $`g(T)`$ is indeterminate. To get the determinate value of $`g`$, we assume in Equation (15) as
$$\frac{k_4}{k}\frac{g_4}{g}=r,$$
(45)
$$\frac{k_4}{k}+\frac{g_4}{g}=s,$$
(46)
where $`r`$ and $`s`$ are constants. From Equations (45) and (46), we have derived
$$k=le^{\frac{(r+s)}{2}T},$$
(47)
$$g=me^{\frac{(rs)}{2}T},$$
(48)
where $`l`$ and $`m`$ are integrating constants. In this case the geometry of the universe (25) reduces to the form
$$ds^2=N^2\left(\frac{1}{X^2}\right)^{\frac{\kappa (\kappa +1)}{(\kappa ^2+2\kappa 1)}}e^{\frac{(sr)}{2}T}(dX^2dT^2+dY^2)$$
$$+\left(\frac{1}{X^2}\right)^{\frac{(\kappa +1)}{(\kappa ^2+2\kappa 1)}}e^{\frac{(r+s)}{2}T}dZ^2,$$
(49)
where
$$N=am\left[\frac{\kappa +1}{(12\kappa \kappa ^2)}\right]^{\frac{\kappa (\kappa +1)}{(\kappa ^2+2\kappa 1)}}and\kappa =\frac{r}{s}.$$
(50)
The effective pressure $`\overline{p}`$ and density $`\rho `$ for the model (49) are given by
$$8\pi \overline{p}=\frac{e^{\frac{(rs)}{2}T}}{N^2X^{\frac{2\kappa (\kappa +1)}{(12\kappa \kappa ^2)}}}\left[\frac{2\kappa ^2(\kappa +1)(\kappa +2)}{(12\kappa \kappa ^2)X^2}\frac{1}{16}(r^2+s^2)\right]\mathrm{\Lambda },$$
(51)
$$8\pi \rho =\frac{e^{\frac{(rs)}{2}T}}{N^2X^{\frac{2\kappa (\kappa +1)}{(12\kappa \kappa ^2)}}}\left[\frac{2\kappa (\kappa +1)(\kappa +2)}{(12\kappa \kappa ^2)X^2}+\frac{1}{8}\{s(2sr)\}\right]+\mathrm{\Lambda }.$$
(52)
On using Equations (4) and (30) in (51), we obtain
$$8\pi (p\xi _0\rho ^n\theta )=\frac{e^{\frac{(rs)}{2}T}}{N^2X^{\frac{2\kappa (\kappa +1)}{(12\kappa \kappa ^2)}}}\left[\frac{2\kappa ^2(\kappa +1)(\kappa +2)}{(12\kappa \kappa ^2)X^2}\frac{1}{16}(r^2+s^2)\right]\mathrm{\Lambda },$$
(53)
where $`\theta `$ is the scalar expansion calculated for the flow vector $`v^i`$ and is given by
$$\theta =\frac{(r+3s)e^{\frac{(rs)}{4}T}}{4\sqrt{N}X^{\frac{\kappa (\kappa +1)}{(12\kappa \kappa ^2)}}}.$$
(54)
### 4.1 Model I: Solution for $`(\xi =\xi _0)`$
In this case, using Eqs. (52), (29) and (54) in Eq. (53), we obtain
$$8\pi (1+\gamma )\rho =\frac{8\pi \xi _0(r+3s)e^{\frac{(rs)}{4}T}}{4\sqrt{N}X^{\frac{\kappa (\kappa +1)}{(12\kappa \kappa ^2)}}}+$$
$$\frac{e^{\frac{(rs)}{2}T}}{N^2X^{\frac{2\kappa (\kappa +1)}{(12\kappa \kappa ^2)}}}\left[\frac{2\kappa (\kappa +1)(\kappa 1)(\kappa +2)}{(12\kappa \kappa ^2)X^2}\frac{1}{16}\left\{r^2+2rs3s^2\right\}\right].$$
(55)
Eliminating $`\rho (t)`$ between (52) and (55), we get
$$(1+\gamma )\mathrm{\Lambda }=\frac{8\pi \xi _0(r+3s)e^{\frac{(rs)}{4}T}}{4\sqrt{N}X^{\frac{\kappa (\kappa +1)}{(12\kappa \kappa ^2)}}}+$$
$$\frac{e^{\frac{(rs)}{2}T}}{N^2X^{\frac{2\kappa (\kappa +1)}{(12\kappa \kappa ^2)}}}\left[\frac{2\kappa (\kappa +1)(\kappa +2)(\kappa +\gamma )}{(12\kappa \kappa ^2)X^2}\frac{1}{16}\left\{r^2+2rs(2+\gamma )s^2(7+4\gamma )\right\}\right].$$
(56)
### 4.2 Model II: Solution for $`(\xi =\xi _0\rho )`$
In this case with the use of Eqs. (52), (29) and (54) in Eq. (53) reduces to
$$8\pi \left[1+\gamma \frac{\xi _0(r+3s)e^{\frac{(rs)}{4}T}}{4\sqrt{N}X^{\frac{\kappa (\kappa +1)}{(12\kappa \kappa ^2)}}}\right]\rho =$$
$$\frac{e^{\frac{(rs)}{2}T}}{N^2X^{\frac{2\kappa (\kappa +1)}{(12\kappa \kappa ^2)}}}\left[\frac{2\kappa (\kappa +1)(\kappa 1)(\kappa +2)}{(12\kappa \kappa ^2)X^2}\frac{1}{16}\left\{r^2+2rs3s^2\right\}\right].$$
(57)
Eliminating $`\rho (t)`$ between (52) and (57), we get
$$[1+\gamma \frac{\xi _0(r+3s)e^{\frac{(rs)}{4}T}}{4\sqrt{N}X^{\frac{\kappa (\kappa +1)}{(12\kappa \kappa ^2)}}}]\mathrm{\Lambda }=\frac{e^{\frac{(rs)}{2}T}}{N^2X^{\frac{2\kappa (\kappa +1)}{(12\kappa \kappa ^2)}}}[\frac{2\kappa ^2(\kappa +1)(\kappa +2)}{(12\kappa \kappa ^2)X^2}$$
$$\frac{1}{16}(r^2+s^2)+\{\frac{2\kappa (\kappa +1)(\kappa +2)}{(12\kappa \kappa ^2)X^2}+\frac{1}{8}(r2s)s\}\times \{\gamma \frac{(r+3s)\xi _0e^{\frac{(rs)}{4}T}}{4\sqrt{N}X^{\frac{\kappa (\kappa +1)}{(12\kappa \kappa ^2)}}}\}].$$
(58)
From Equations (56) and (58), we observe that cosmological constant $`\mathrm{\Lambda }`$ may be positive or negative under specific conditions. A negative cosmological constant adds to the attractive gravity of matter, therefore, universe with a negative cosmological constant are invariably doomed to recollapse. A positive cosmological constant resists the attractive gravity of matter due to its negative pressure. For most universe, the positive cosmological constant eventually dominates over the attraction of matter and drives the universe to expands exponentially.
### 4.3 Some physical aspects of the models
The coefficient of shear $`\sigma `$, non-vanishing physical components of conformal curvature tensor $`C_{(ijkl)}`$, non-vanishing component of electromagnetic field tensor $`F_{ij}`$ and component of charge density $`J^2`$ for the model (49) are given as:
$$\sigma ^2=\frac{r^2e^{\frac{(rs)}{2}T}}{8NX^{\frac{2\kappa (\kappa +1)}{(12\kappa \kappa ^2)}}},$$
(59)
$$C_{(1212)}=C_{(3434)}=\frac{e^{\frac{(rs)}{2}T}}{6NX^{\frac{2\kappa (\kappa +1)}{(12\kappa \kappa ^2)}}}\left[\frac{2\kappa (\kappa +1)(\kappa 1)}{(12\kappa \kappa ^2)^2X^2}\frac{r^2}{2}\right],$$
(60)
$$C_{(1313)}=C_{(2424)}=\frac{e^{\frac{(rs)}{2}T}}{6NX^{\frac{2\kappa (\kappa +1)}{(12\kappa \kappa ^2)}}}\left[\frac{4\kappa (\kappa +1)(\kappa 1)}{(12\kappa \kappa ^2)^2X^2}+\frac{r^2}{4}\right],$$
(61)
$$C_{(1224)}=C_{(1334)}=\frac{re^{\frac{(rs)}{2}T}}{4NX^{\frac{2\kappa (\kappa +1)}{(12\kappa \kappa ^2)}}},$$
(62)
$$C_{(2323)}=C_{(1414)}=\frac{e^{\frac{(rs)}{2}T}}{3NX^{\frac{2\kappa (\kappa +1)}{(12\kappa \kappa ^2)}}}\left[\frac{\kappa (\kappa +1)^2(\kappa 1)}{(12\kappa \kappa ^2)^2X^2}+\frac{1}{8}\left\{r(rs)\right\}\right],$$
(63)
$$F_{12}=\frac{1}{\sqrt{8\pi }}\left[\frac{2\kappa (\kappa +1)(1\kappa )(\kappa +2)}{(12\kappa \kappa ^2)X^2}\frac{rs}{4}\right]^{\frac{1}{2}}B,$$
(64)
$$J^2=\frac{e^{\frac{(rs)}{2}T}}{N^2X^{\frac{(\kappa +1)(3\kappa +1)}{(12\kappa \kappa ^2)}}}\times $$
$$\left[\frac{(1\kappa ^2)}{(12\kappa \kappa ^2)}\psi ^{\frac{1}{2}}X^{\frac{2\kappa }{(12\kappa \kappa ^2)}}+\frac{2\kappa (\kappa +1)(\kappa 1)(\kappa +2)}{(12\kappa \kappa ^2)^2}\psi ^{\frac{1}{2}}X^{\frac{2(\kappa ^2+3\kappa 1)}{(12\kappa \kappa ^2)}}\right],$$
(65)
where
$$\psi =\left[\frac{2\kappa (\kappa +1)(1\kappa )(\kappa +2)}{(12\kappa \kappa ^2)X^2}\frac{rs}{4}\right].$$
(66)
The rotation $`\omega `$ is identically zero. The models start expanding at $`T=0`$ and continue till $`T=\mathrm{}`$ for $`N>0`$. The expansion stops when $`T=\mathrm{}`$ or $`r=3s`$. The models represent shearing, non-rotating and Petrov type I non-degenerate in general.
## 5 Conclusions
We have obtained a new class of plane-symmetric inhomogeneous cosmological models of electromagnetic bulk viscous fluid as the source of matter. Generally the models represent expanding, shearing, non-rotating and Petrov type I non-degenerate universe in which the flow vector is geodetic. In all these models, we observe that they do not approach isotropy for large values of time.
The cosmological constants in all models given in Sections $`3`$ are decreasing functions of time and they all approach a small value at late times. The values of cosmological “constant” for these models are found to be small and positive which are supported by the results from recent supernova Ia observations recently obtained by the High-Z Supernova Team and Supernova Cosmological Project ( Garnavich et al.; Perlmutter et al.; Riess et al.; Schmidt et al. ).
## Acknowledgements
One of the authors (A. Pradhan) would like to thank the Inter-University Centre for Astronomy and Astrophysics, Pune, India for providing facility under Associateship Programme where part of this work was carried out. Authors thank the anonymous referee for many helpful comments which helped in improving the presentation of this paper. Authors also thank S. K. Srivastava and Saibal Ray for useful discussions. |
warning/0506/cs0506094.html | ar5iv | text | # Universal Codes as a Basis for Nonparametric Testing of Serial Independence for Time Series
## I Introduction
Nonparametric testing for independence of time series is very important in statistical applications. There is an extensive literature dealing with nonparametric independence testing; a quite full review can be found in .
In this paper, we consider a source (or process), which generates elements from a finite set $`A`$ and two following hypotheses: $`H_0`$ is that the source is Markovian one, which memory (or connectivity) not larger than $`m,(m0),`$ and the alternative hypothesis $`H_1`$ that the sequence is generated by a stationary and ergodic source, which differs from the source under $`H_0`$. The testing should be based on a sample $`x_1\mathrm{}x_t`$ generated by the source.
For example, the sequence $`x_1\mathrm{}x_t`$ might be a DNA-string and one can consider the question about the depth of the statistical dependence.
We suggest a family of tests that are based on so called universal predictors (or universal data compression methods). The Type I error of the suggested tests is not larger than a given $`\alpha (\alpha (0,1))`$ for any source under $`H_0,`$ whereas the Type II error for any source under $`H_1`$ tends to 0, when the sample size $`t`$ grows.
The suggested tests are based on results and ideas of Information Theory and, especially, those of the universal coding theory. Informally, the main idea of the tests can be described as follows. Suppose that the source generates letters from an alphabet $`A`$ and one wants to test $`H_0`$ (the source is Morkovian with memory $`m,m0.`$ ) First we recall that there exist universal codes which, informally speaking can ”compress” any sequence generated by a stationary and ergodic source, to the length $`th_{\mathrm{}}`$ bits, where $`h_{\mathrm{}}`$ is the limit Shannon entropy and $`t`$ tends to infinity. Second, it is well known in information theory that $`h_{\mathrm{}}`$ equals $`m`$th-order (conditional) Shannon entropy $`h_m`$, if $`H_0`$ is true, and $`h_{\mathrm{}}`$ is strictly less than $`h_m`$ if $`H_1`$ is true. So, the following test looks like natural: Compress the sample sequence $`x_1\mathrm{}x_t`$ by a universal code and compare the lengths of the obtained file with $`th_m^{},`$ where $`h_m^{}`$ is an estimate of $`h_m.`$ If the length of the compressed file is significantly less than $`th_m^{},`$ then the hypothesis $`H_0`$ should be rejected.
This is no surprise that the results and ideas of a universal coding theory can be applied to some classical problems of mathematical statistics. In fact, methods of universal coding (and a closely connected universal prediction) are intended to extract information from observed data in order to compress (or predict) data efficiently in a case where the source statistics is unknown. Recently such a connection between universal coding and mathematical statistics was used in for estimating the order of Markov sources and for constructing efficient tests for randomness, i.e. for testing the hypothesis $`\widehat{H}_0`$ that a sequence is generated by a Bernoulli source and all letters have equal probabilities against $`\widehat{H}_1`$ that the sequence is generated by a stationary and ergodic source, which differs from the source under $`\widehat{H}_0,`$ see .
The outline of the paper is as follows. The next part contains definitions and necessary information from the theory of universal coding and universal prediction. Part three is devoted to the testing of the above described hypotheses. All proofs are given in the appendix.
## II Definitions and Preliminaries.
Consider an alphabet $`A=\{a_1,\mathrm{},a_n\}`$ with $`n2`$ letters and denote by $`A^t`$ the set of words $`x_1\mathrm{}x_t`$ of length $`t`$ from $`A`$. Let $`\mu `$ be a source which generates letters from $`A`$. Formally, $`\mu `$ is a probability distribution on the set of words of infinite length or, more simply, $`\mu =(\mu ^t)_{t1}`$ is a consistent set of probabilities over the sets $`A^t;t1`$. By $`M_{\mathrm{}}(A)`$ we denote the set of all stationary and ergodic sources, which generate letters from $`A.`$ Let $`M_k(A)M_{\mathrm{}}(A)`$ be the set of Markov sources with memory (or connectivity) $`k,k0.`$ More precisely, by definition $`\mu M_k(A)`$ if
$$\mu (x_{t+1}=a_{i_1}/x_t=a_{i_2},x_{t1}=a_{i_3},\mathrm{},x_{tk+1}=a_{i_{k+1}},\mathrm{})$$
$$=\mu (x_{t+1}=a_{i_1}/x_t=a_{i_2},x_{t1}=a_{i_3},\mathrm{},x_{tk+1}=a_{i_{k+1}})$$
for all $`tk`$ and $`a_{i_1},a_{i_2},\mathrm{}A.`$ By definition, $`M_0(A)`$ is the set of all Bernoulli (or i.i.d.) sources over $`A`$.
2.1 Universal prediction.
Now we briefly describe results and methods of universal coding and prediction, which will be used later. Let a source generate a message $`x_1\mathrm{}x_{t1}x_t\mathrm{}`$ and let $`\nu ^t(a)`$ denote the count of letter $`a`$ occurring in the word $`x_1\mathrm{}x_{t1}x_t`$. After the first $`t`$ letters $`x_1,\mathrm{},x_{t1},x_t`$ have been processed the following letter $`x_{t+1}`$ needs to be predicted. By definition, a prediction is a set of non-negative numbers $`\pi (a_1|x_1\mathrm{}x_t),\mathrm{},\pi (a_n|x_1\mathrm{}x_t)`$ which are estimates of the unknown conditional probabilities $`p(a_1|x_1\mathrm{}x_t),\mathrm{},p(a_n|x_1\mathrm{}x_t)`$, i.e. of the probabilities $`p(x_{t+1}=a_i|x_1\mathrm{}x_t)`$; $`i=1,\mathrm{},n`$.
Laplace suggested the following predictor:
$$L(a|x_1\mathrm{}x_t)=(\nu ^t(a)+1)/(t+|A|),$$
(1)
where $`|A|`$ is the number of letters in the alphabet $`A,`$ see . For example, if $`A=\{0,1\},x_1\mathrm{}x_5=01010,`$ then the Laplace prediction is as follows: $`L(x_6=0|01010)=(3+1)/(5+2)=4/7,L_0(x_6=1|01010)=(2+1)/(5+2)=3/7.`$
In Information Theory the error of prediction often is estimated by the the Kullback-Leibler (K-L) divergence between a distribution $`p`$ and its estimation. Consider a source $`p`$ and a predictor $`\gamma `$. The error is characterized by the divergence
$$\rho _{\gamma ,p}(x_1\mathrm{}x_t)=\underset{aA}{}p(a|x_1\mathrm{}x_t)\mathrm{log}\frac{p(a|x_1\mathrm{}x_t)}{\gamma (a|x_1\mathrm{}x_t)}.$$
(2)
(Here and below $`\mathrm{log}\mathrm{log}_2`$.) It is well known that for any distributions $`p`$ and $`\gamma `$ the K-L divergence is nonnegative and equals 0 if and only if $`p(a)=\gamma (a)`$ for all $`a,`$ see, for ex., , that is why the K-L divergence is a natural estimate of the prediction error. For fixed $`t`$, $`\rho _{\gamma ,p}`$ is a random variable, because $`x_1,x_2,\mathrm{},x_t`$ are random variables. We define the average error at time $`t`$ by
$$\rho ^t(p\gamma )=E\left(\rho _{\gamma ,p}()\right)=\underset{x_1\mathrm{}x_tA^t}{}p(x_1\mathrm{}x_t)\rho _{\gamma ,p}(x_1\mathrm{}x_t).$$
It is known that the error of Laplace predictor goes to 0 for any Bernoulli source $`p`$. More precisely, it is proven that
$$\rho ^t(pL)<(|A|1)/(t+1)$$
(3)
for any source $`p;`$ .
Obviously, the convergence to 0 of a predictor’s error for any source from some set $`M`$ is an important property. For example, we can see from (3) that it is true for the Laplace predictor and the set of Bernoulli sources $`M_0(A)`$. Unfortunately, it is known that a predictor, which error (2) goes to 0 for any stationary and ergodic source, does not exist. More precisely, for any predictor $`\gamma `$ there exists such a stationary and ergodic source $`\stackrel{~}{p},`$ that $`lim_t\mathrm{}sup(\rho _{\gamma ,\stackrel{~}{p}}(x_1\mathrm{}x_t)const>0`$ with probability 1; . (See also , where this result is generalized and a history of its discovery is described. In particular, they found out that such a result was described by Bailey in his unpublished thesis). That is why it is difficult to use (2) for comparison of different predictors. On the other hand, it is shown in that there exists such a predictor $`R`$, that the following average $`t^1_{i=1}^t\rho _{R,p}(x_1\mathrm{}x_t)`$ goes to 0 (with probability 1 ) for any stationary and ergodic source $`p,`$ where $`t`$ goes to infinity. That is why we will focus our attention on such averages. First, we define for any predictor $`\pi `$ the following probability distribution
$$\pi (x_1\mathrm{}x_t)=\underset{i=1}{\overset{t}{}}\pi (x_i|x_1\mathrm{}x_{i1}).$$
For example, we obtain for the Laplace predictor $`L`$ that $`L(0101)=\frac{1}{2}\frac{1}{3}\frac{1}{2}\frac{2}{5}=\frac{1}{30},`$ see (1). Then, by analogy with (2) we will estimate the error by K-L divergence and define
$$\rho _{\gamma ,p}(x_1\mathrm{}x_t)=t^1(\mathrm{log}(p(x_1\mathrm{}x_t)/\gamma (x_1\mathrm{}x_t))$$
(4)
and
$$\overline{\rho }_t(\gamma ,p)=t^1\underset{x_1\mathrm{}x_tA^t}{}p(x_1\mathrm{}x_t)\mathrm{log}(p(x_1\mathrm{}x_t)/\gamma (x_1\mathrm{}x_t)).$$
(5)
For example, from those definitions and (3) we obtain the following estimation for Laplace predictor $`L`$ and any Bernoulli source: $`\overline{\rho }_t(L,p)<((|A|1)\mathrm{log}t+c)/t,`$ where $`c`$ is a constant.
The universal predictors will play a key rule in suggested below tests. By definition, a predictor $`\gamma `$ is called a universal (in average) for a class of sources $`M`$, if for any $`pM`$ the error $`\overline{\rho }_t(\gamma ,p)`$ goes to 0, where $`t`$ goes to infinity. A predictor $`\gamma `$ is called universal with probability 1, if the error $`\rho _{\gamma ,p}(x_1\mathrm{}x_t)`$ goes to 0 not only in average, but for almost all sequences $`x_1x_2\mathrm{}`$. For short, we will say that the predictor (or probability distribution) $`\gamma `$ is universal, if $`lim_t\mathrm{}\rho _{\gamma ,p}(x_1\mathrm{}x_t)=0`$ is valid with probability 1 for any stationary and ergodic source (i.e. for any $`pM_{\mathrm{}}(A)`$). Now there are quite many known universal predictors. One of the first such predictors is described in .
2.1 Universal coding.
This short subparagraph is intended to give some explanation about why and how methods of data compression can be used for testing of independence. The point is that the prediction problem is deeply connected with the theory of universal coding. Moreover, practically used data compression methods (or so-called archivers) can be directly applied for testing.
Let us give some definitions. Let, as before, $`A`$ be a finite alphabet and, by definition, $`A^{}=_{n=1}^{\mathrm{}}A^n`$ and $`A^{\mathrm{}}`$ is the set of all infinite words $`x_1x_2\mathrm{}`$ over the alphabet $`A`$. A data compression method (or code) $`\phi `$ is defined as a set of mappings $`\phi _n`$ such that $`\phi _n:A^n\{0,1\}^{},n=1,2,\mathrm{}`$ and for each pair of different words $`x,yA^n`$ $`\phi _n(x)\phi _n(y).`$ Informally, it means that the code $`\phi `$ can be applied for compression of each message of any length $`n,n>0`$ over alphabet $`A`$ and the message can be decoded if its code is known. One more restriction is required in Information Theory. Namely, it is required that each sequence $`\phi _n(x_1)\phi _n(x_2)\mathrm{}\phi _n(x_r),r1,`$ of encoded words from the set $`A^n,n1,`$ can be uniquely decoded into $`x_1x_2\mathrm{}x_r`$. Such codes are called uniquely decodable. For example, let $`A=\{a,b\}`$, the code $`\psi _1(a)=0,\psi _1(b)=00,`$ obviously, is not uniquely decodable. (Indeed, the word $`000`$ can be decoded in both $`ab`$ and $`ba.`$) It is well known that if a code $`\phi `$ is uniquely decodable then the lengths of the codewords satisfy the following inequality (the Kraft inequality):
$$\mathrm{\Sigma }_{uA^n}\mathrm{\hspace{0.25em}2}^{|\phi _n(u)|}1,$$
see, for ex., . It will be convenient to reformulate this property as follows:
Claim 1. Let $`\phi `$ be a uniquely decodable code over an alphabet $`A`$. Then for any integer $`n`$ there exists a measure $`\mu _\phi `$ on $`A^n`$ such that
$$\mathrm{log}\mu _\phi (u)|\phi (u)|$$
(6)
for any $`u`$ from $`A^n.`$ (Obviously, it is true for the measure $`\mu _\phi (u)=2^{|\phi (u)|}/\mathrm{\Sigma }_{uA^n}\mathrm{\hspace{0.25em}2}^{|\phi (u)|}.`$) It is known in Information Theory that sequences $`x_1\mathrm{}x_t,`$ generated by a stationary and ergodic source $`p,`$ can be ”compressed” till the length $`\mathrm{log}p(x_1\mathrm{}x_t)`$ bits. There exist so-called universal codes, which, in a certain sense, are the best ”compressors” for all stationary and ergodic sources. The formal definition is as follows: A code $`\phi `$ is universal if for any stationary and ergodic source $`p`$
$$\underset{t\mathrm{}}{lim}t^1(\mathrm{log}p(x_1\mathrm{}x_t)|\phi (x_1\mathrm{}x_t)|=\mathrm{\hspace{0.17em}0}$$
with probability 1. So, informally speaking, the universal codes estimate the probability characteristics of the source $`p`$ and use them for efficient ”compression”.
## III The Tests.
In this paragraph we describe the suggested tests. First, we give some definitions. Let $`v`$ be a word $`v=v_1\mathrm{}v_k,kt,v_iA.`$ Denote the rate of a word $`v`$ occurring in the sequence $`x_1x_2\mathrm{}x_k`$ , $`x_2x_3\mathrm{}x_{k+1}`$, $`x_3x_4\mathrm{}x_{k+2}`$, $`\mathrm{}`$, $`x_{tk+1}\mathrm{}x_t`$ as $`\nu ^t(v)`$. For example, if $`x_1\mathrm{}x_t=000100`$ and $`v=00,`$ then $`\nu ^6(00)=3`$. Now we define for any $`k0`$ a so-called empirical Shannon entropy of order $`k`$ as follows:
$$h_k^{}(x_1\mathrm{}x_t)=$$
(7)
$$\frac{1}{(tk)}\underset{vA^k}{}\overline{\nu }^t(v)\underset{aA}{}(\nu ^t(va)/\overline{\nu }^t(v))\mathrm{log}(\nu ^t(va)/\overline{\nu }^t(v)),$$
where $`k<t`$ and $`\overline{\nu }^t(v)=_{aA}\nu ^t(va).`$ In particular, if $`k=0`$, we obtain
$$h_0^{}(x_1\mathrm{}x_t)=\frac{1}{t}\underset{aA}{}\nu ^t(a)\mathrm{log}(\nu ^t(a)/t),$$
The suggested test is as follows.
*Let $`\sigma `$ be any probability distribution over $`A^t.`$ By definition, the hypothesis $`H_0`$ is accepted if*
$$(tm)h_m^{}(x_1\mathrm{}x_t)\mathrm{log}(1/\sigma (x_1\mathrm{}x_t))\mathrm{log}(1/\alpha ),$$
(8)
*where $`\alpha (0,1).`$ Otherwise, $`H_0`$ is rejected.* We denote this test by $`\mathrm{{\rm Y}}_{\alpha ,\sigma ,m}^t.`$
Theorem. i) For any probability distribution (or predictor) $`\sigma `$ the Type I error of the test $`\mathrm{{\rm Y}}_{\alpha ,\sigma ,m}^t`$ is less than or equal to $`\alpha ,\alpha (0,1)`$.
ii) If $`\sigma `$ is a universal predictor (measure) (i.e., by definition, for any $`pM_{\mathrm{}}(A)`$
$$\underset{t\mathrm{}}{lim}t^1(\mathrm{log}p(x_1\mathrm{}x_t)\mathrm{log}(1/\sigma (x_1\mathrm{}x_t))=\mathrm{\hspace{0.17em}0}$$
(9)
with probability 1), then the Type II error goes to 0, where $`t`$ goes to infinity.
The proof is given in Appendix.
Comment. Let $`\phi `$ be a uniquely decodable code (or a data compression method). Define the test $`\widehat{\mathrm{{\rm Y}}}_{\alpha ,\phi ,m}^t`$ as follows: The hypothesis $`H_0`$ is accepted if
$$(tm)h_m^{}(x_1\mathrm{}x_t)|\phi (x_1\mathrm{}x_t)|\mathrm{log}(1/\alpha ),$$
(10)
where $`\alpha (0,1).`$ Otherwise, $`H_0`$ is rejected.
We immediately obtain from the theorem 1 and the claim 1 the following statement.
Claim 2. i) For any uniquely decodable code $`\phi `$ the Type I error of the test $`\widehat{\mathrm{{\rm Y}}}_{\alpha ,\phi ,m}^t`$ is less than or equal to $`\alpha ,\alpha (0,1)`$.
ii) If $`\phi `$ is a universal code, then the Type II error goes to 0, where $`t`$ goes to infinity.
## IV Conclusion.
The described above tests can be based on known universal codes (or so-called archivers) which are used for text compression everywhere. It is important to note that, on the one hand, the universal codes and archivers are based on results of Information Theory, the theory of algorithms and some other branches of mathematics; see, for example, probability . On the other hand, the archivers have shown high efficiency in practise as compressors of texts, DNA sequences and many other types of real data. In fact, the archivers can find many kinds of latent regularities, that is why they look like a promising tool for independence testing and its generalizations.
The natural question is a possibility of generalization of the suggested tests for a case of an infinite source alphabet $`A`$ (say, $`A`$ is a metric space.) Apparently, such a generalization can be done for a case of independence testing, if we will use known methods of partitioning; . But we do not know how to generalize the suggested tests for a case where $`H_0`$ is that the source is Markovian. The point is that the partitioning can increase the source memory. For example, even if the alphabet $`A`$ contains three letters and we combine two of them in one subset (i.e. a new letter) the memory of the obtained source can increase till infinity. Hence, the generalization to Markov sources with infinite alphabet can be considered as an open problem.
## V Appendix.
Proof of Theorem. First we show that for any Bernoulli source $`\tau ^{}`$ and any word $`x_1\mathrm{}x_tA^t,t>1,`$ the following inequality is valid:
$$\tau ^{}(x_1\mathrm{}x_t)=\underset{aA}{}\tau (a)^{\nu ^t(a)}\underset{aA}{}(\nu ^t(a)/t)^{\nu ^t(a)}$$
(11)
Indeed, the equality is true, because $`\tau ^{}`$ is a Bernoulli measure. The inequality follows from the well known inequality $`_{aA}p(a)\mathrm{log}(p(a)/q(a))0,`$ for K-L divergence, which is true for any distributions $`p`$ and $`q`$ (see, for ex., ). So, if $`p(a)=\nu ^t(a)/t`$ and $`q(a)=\tau ^{}(a),`$ then
$$\underset{aA}{}\frac{\nu ^t(a)}{t}\mathrm{log}\frac{(\nu ^t(a)/t)}{\tau (a)}0.$$
From the last inequality we obtain (11).
Let now $`\tau `$ belong to $`M_m(A),m>0.`$ We will prove that for any $`x_1\mathrm{}x_t`$
$$\tau (x_1\mathrm{}x_t)\underset{uA^m}{}\underset{aA}{}(\nu ^t(ua)/\overline{\nu }^t(u))^{\nu ^t(ua)}.$$
(12)
Indeed, we can present $`\tau (x_1\mathrm{}x_t)`$ as
$$\tau (x_1\mathrm{}x_t)=\tau _{\mathrm{}}(x_1\mathrm{}x_m)\underset{uA^m}{}\underset{aA}{}\tau (a/u)^{\nu ^t(ua)},$$
where $`\tau _{\mathrm{}}(x_1\mathrm{}x_m)`$ is the limit probability of the word $`x_1\mathrm{}x_m.`$ From the last equality we can see that
$$\tau (x_1\mathrm{}x_t)\underset{uA^m}{}\underset{aA}{}\tau (a/u)^{\nu ^t(ua)}.$$
Taking into account the inequality (11), we obtain
$$\underset{aA}{}\tau (a/u)^{\nu ^t(ua)}\underset{aA}{}(\nu ^t(ua)/\overline{\nu }^t(u))^{\nu ^t(ua)}$$
for any word $`u`$. So, from the last two inequalities we obtain (12).
It will be convenient to define an auxiliary measure on $`A^t`$ as follows:
$$\pi _m(x_1\mathrm{}x_t)=\mathrm{\Delta }\mathrm{\hspace{0.25em}2}^{th_m^{}(x_1\mathrm{}x_t)},$$
(13)
where $`x_1\mathrm{}x_tA^t`$ and
$$\mathrm{\Delta }=(\underset{x_1\mathrm{}x_tA^t}{}\mathrm{\hspace{0.17em}2}^{th_m^{}(x_1\mathrm{}x_t)})^1.$$
If we take into account that
$$2^{(tm)h_m^{}(x_1\mathrm{}x_t)}=\underset{uA^m}{}\underset{aA}{}(\nu ^t(ua)/\overline{\nu }^t(u))^{\nu ^t(ua)},$$
we can see from (12) and (13) that, for any measure $`\tau M_m(A)`$ and any $`x_1\mathrm{}x_tA^t,`$
$$\tau (x_1\mathrm{}x_t)\pi _m(x_1\mathrm{}x_t)/\mathrm{\Delta }.$$
(14)
Let us denote the critical set of the test $`\mathrm{{\rm Y}}_{\alpha ,\sigma ,m}^t`$ as $`C_\alpha `$ i.e., by definition,
$$C_\alpha =\{x_1\mathrm{}x_t:(tm)h_m^{}(x_1\mathrm{}x_t)\mathrm{log}(1/\sigma (x_1\mathrm{}x_t))$$
$$>\mathrm{log}(1/\alpha )\}.$$
(15)
From (14) and this definition we can see that for any measure $`\tau M_m(A)`$
$$\tau (C_\alpha )\pi _m(C_\alpha )/\mathrm{\Delta }.$$
(16)
From the definitions (15) and (13) we obtain
$$C_\alpha =\{x_1\mathrm{}x_t:\mathrm{\hspace{0.33em}2}^{(tm)h_m^{}(x_1\mathrm{}x_t)}>(\alpha \sigma (x_1\mathrm{}x_t))^1\}$$
$$=\{x_1\mathrm{}x_t:(\pi _m(x_1\mathrm{}x_t)/\mathrm{\Delta })^1>(\alpha \sigma (x_1\mathrm{}x_t))^1\}.$$
Finally,
$$C_\alpha =\{x_1\mathrm{}x_t:\sigma (x_1\mathrm{}x_t)>\pi _m(x_1\mathrm{}x_t)/(\alpha \mathrm{\Delta })\}.$$
(17)
The following chain of inequalities and equalities is valid:
$$1\underset{x_1\mathrm{}x_tC_\alpha }{}\sigma (x_1\mathrm{}x_t)\underset{x_1\mathrm{}x_tC_\alpha }{}\pi _m(x_1\mathrm{}x_t)/(\alpha \mathrm{\Delta })$$
$$=\pi _m(C_\alpha )/(\alpha \mathrm{\Delta })\tau (C_\alpha )\mathrm{\Delta }/(\alpha \mathrm{\Delta })=\tau (C_\alpha )/\alpha .$$
(Here both equalities and the first inequality are obvious, the second inequality and the third one follow from (17) and (16), correspondingly.) So, we obtain that $`\tau (C_\alpha )\alpha `$ for any measure $`\tau M_m(A).`$ Taking into account that $`C_\alpha `$ is the critical set of the test, we can see that the probability of the Type I error is not greater than $`\alpha .`$ The first claim of the theorem is proven.
The proof of the second statement of the theorem will be based on some results of Information Theory. The $`t`$ order conditional Shannon entropy is defined as follows:
$$h_t(p)=\underset{x_1\mathrm{}x_tA^t}{}p(x_1\mathrm{}x_t)$$
$$\underset{aA}{}p(a/x_1\mathrm{}x_t)\mathrm{log}p(a/x_1\mathrm{}x_t),$$
(18)
where $`pM_{\mathrm{}}(A).`$ It is known that for any $`pM_{\mathrm{}}(A)`$ firstly, $`\mathrm{log}|A|h_0(p)h_1(p)\mathrm{},`$ secondly, there exists the following limit Shannon entropy $`h_{\mathrm{}}(p)=lim_t\mathrm{}h_t(p)`$, thirdly, $`lim_t\mathrm{}t^1\mathrm{log}p(x_1\mathrm{}x_t)=h_{\mathrm{}}(p)`$ with the probability 1 and, finally, $`h_m(p)`$ is strictly greater than $`h_{\mathrm{}}(p),`$ if the memory of $`p`$ is larger $`m`$, (i.e. $`pM_{\mathrm{}}(A)M_m(A)`$), see, for example, .
Taking into account the definition of the universal predictor (see (9)), we obtain from the above described properties of the entropy that
$$\underset{t\mathrm{}}{lim}t^1\mathrm{log}\sigma (x_1\mathrm{}x_t)=h_{\mathrm{}}(p)$$
(19)
with probability 1. It can be seen that $`h_m^{}`$ (10) is a consistent estimate for the $`m`$order Shannon entropy (18), i.e. $`lim_t\mathrm{}h_m^{}(x_1\mathrm{}x_t)=h_m(p)`$ with probability 1; see . Having taken into account that $`h_m(p)>h_{\mathrm{}}(p)`$ and (19) we obtain from the last equality that $`lim_t\mathrm{}((tm)h_m^{}(x_1\mathrm{}x_t)\mathrm{log}(1/\sigma (x_1\mathrm{}x_t)))=\mathrm{}.`$ This proves the second statement of the theorem. |
warning/0506/cond-mat0506782.html | ar5iv | text | # Integrable dynamics of coupled Fermi–Bose condensates
## I Introduction
A number of closely related non-stationary problems have come up recently in different contexts. In these problems the goal is to describe the dynamics of a many-body system following a sudden perturbation that drove the system out of an equilibrium. The system in question can be a BCS superconductorander1 ; galaiko ; volkov ; galperin ; shumeiko ; barankov1 ; simons ; amin ; BCS ; shortBCS ; leggett , coupled Fermi-Bose condensatesandreev ; barankov , or a single electronic spin interacting with many nuclear spins (the central spin model)stamp ; khaetskii ; schliemann ; nazarov1 ; coish ; koka ; lukin . A common feature of all these problems is that they can be formulated in terms of spin Hamiltonians, which belong to a class of integrable systems known as Gaudin magnetsgaudin ; sklyanin ; vadim . It turns out that this fact enables one to solve for the evolution of all these systems using the same approach.
In the present paper we focus on a particular example of such a problem. Namely, we study the dynamics of a fermionic condensate interacting with a single bosonic mode. Recent studies of this problemandreev ; barankov were motivated by experimentsregal ; zwierlein1 ; kinast ; zwierlein2 on cold fermion pairing. The idea was that near the Feshbach resonance feshbach1 ; feshbach2 these systems can be modeled by a fermionic condensate of atoms strongly coupled to the lowest energy bosonic mode.
We assume that the system has been prepared in a nonequilibrium state at $`t0`$ and study the subsequent evolution for $`t>0`$. Our main result is a complete solution for the dynamics of the system. We also derive a full set of its integrals of motion. It turns out that the mean-field time evolution of the coupled Fermi-Bose condensates can be mapped onto the corresponding BCS problem, which was solved previously in Refs. BCS, ; shortBCS, . For this reason, the emphasis in the present paper is on this mapping and detailed or alternative derivations that have been largely omitted in Refs. BCS, ; shortBCS, .
The fermion-boson condensate is described by the following Hamiltonian:
$$\widehat{H}=\underset{j,\sigma }{}ϵ_j\widehat{c}_{j\sigma }^{}\widehat{c}_{j\sigma }+\omega \widehat{b}^{}\widehat{b}+g\underset{j}{}\left(\widehat{b}^{}\widehat{c}_j\widehat{c}_j+\widehat{b}\widehat{c}_j^{}\widehat{c}_j^{}\right),$$
(1)
where $`ϵ_j`$ are the single-particle energy levels and the operators $`\widehat{c}_{j\sigma }^{}`$ ($`\widehat{c}_{j\sigma }`$) create (annihilate) a fermion of one of the two species $`\sigma =`$ or $``$ in an orbital eigenstate of energy $`ϵ_j`$. Eigenstates $`|j`$ and $`|j`$ are related by the time-reversal symmetryander . For example, if the single-particle potential is translationally invariant, $`|j=|𝐩`$ and $`|j=|𝐩`$. Operators $`\widehat{b}^{}`$ ($`\widehat{b}`$) create (annihilate) quanta of the bosonic field.
We study the dynamics of the fermion-oscillator model (1) in the mean-field approximation. This amounts to treating the bosonic field classically, i.e. replacing operators $`\widehat{b}^{}`$ and $`\widehat{b}`$ with $`c`$-numbers in the Heisenberg equations of motion for Hamiltonian (1). This procedure is expected to be exact as long as the bosonic mode is macroscopically populated. It turns out that in this approximation the dynamics coincides with that of a classical Hamiltonian system, which can be mapped onto the corresponding BCS problem (see below). The classical dynamical variables are the time-dependent quantum-mechanical expectation values $`\widehat{c}_j\widehat{c}_j`$, $`\widehat{b}`$, and $`_\sigma \widehat{c}_{j\sigma }^{}\widehat{c}_{j\sigma }`$. We will see below that the mean-field approximation is also equivalent to the usual way of obtaining the classical limit – by replacing operators with classical variables and their commutators with Poisson brackets.
Here we assume that the number of energy levels, $`n`$, interacting with the bosonic field in the system (1) is arbitrary large, but finite. In this case the typical evolution is quasi-periodic with $`n+1`$ incommensurate frequencies. The system exhibits irregular multi-frequency oscillations ergodically exploring the part of the phase-space allowed by the conservation laws, returning arbitrarily close to its initial state at irregular time intervals. For finite $`n`$, the solution, albeit explicit, is rather complicated. However, it considerably simplifies in the thermodynamic limit $`n\mathrm{}`$. In this limit, the return time diverges for most physical initial conditions, while the solution decays to a simple limiting dynamics at large timesunpub . For example, following an initial decay, $`|\widehat{b}(t)|`$ asymptotes either to a constant value (cf. Refs. galaiko, ; volkov, ; simons, ; leggett, ) or to an oscillatory behavior characterized by only a few frequencies (cf. Fig. 2 in Refs. galperin, ; barankov1, ).
The fermion-oscillator model (1) can be viewed as a generalization of the Dicke (Tavis-Cummings) model of Quantum Optics. The latter can be obtained from (1) in the zero fermionic bandwidth limit, i.e. when all single-particle levels $`ϵ_j`$ are degenerate, $`ϵ_j=\mu `$. To see this, it is useful to reformulate the model (1) as a spin-oscillator model using Anderson’s pseudospin representationander1 . Pseudospins are defined as
$$2\widehat{K}_j^z=\widehat{n}_j1,\widehat{K}_j^{}=\widehat{c}_j\widehat{c}_j,\widehat{K}_j^+=\widehat{c}_j^{}\widehat{c}_j^{},$$
(2)
where $`\widehat{n}_j=_\sigma \widehat{c}_{j\sigma }^{}\widehat{c}_{j\sigma }`$. Operators $`\widehat{𝐊}_j`$ have all properties of spin-1/2 on the subspace of unoccupied and doubly occupied (unblocked) levels $`ϵ_j`$. Singly occupied (blocked) levels do not interact with the bosonic field and are decoupled from the dynamics. In terms of pseudospins the Hamiltonian (1) for $`n`$ unblocked levels takes the form
$$\widehat{H}=\underset{j=0}{\overset{n1}{}}2ϵ_j\widehat{K}_j^z+\omega \widehat{b}^{}\widehat{b}+g\underset{j=0}{\overset{n1}{}}\left(\widehat{b}^{}\widehat{K}_j^{}+\widehat{b}\widehat{K}_j^+\right).$$
(3)
In the zero bandwidth limit the Hamiltonian (3) describes an interaction of a single collective spin $`\widehat{𝐓}=_j\widehat{𝐊}_j`$ with a harmonic oscillator
$$\widehat{H}_{Dicke}=2\mu \widehat{T}_z+\omega \widehat{b}^{}\widehat{b}+g(\widehat{b}^{}\widehat{T}_{}+\widehat{b}\widehat{T}_+).$$
(4)
This model was suggested by R. H. Dicke in 1953 dicke . Its spectrum was obtained exactly by M. Tavis and F. W. Cummings in 1967 TC . In the mean-field approximation it becomes a classical model for the dynamical variables $`\widehat{𝐓}(t)`$ and $`\widehat{b}(t)`$ with only two degrees of freedom (see also below). The mean-field solution for the time evolution was outlined by Dickedicke and derived in detail by R. Bonifacio and G. Preparata in 1969 BP . In fact, the resulting classical problem is that of a spherical pendulum. The solution for $`\widehat{𝐓}(t)`$ and $`\widehat{b}(t)`$ is in terms of elliptic functions and (apart from the azimuthal motion of the pendulum) is simply periodic BP .
Remarkably, a more general many-body model (1,3) also turns out to be integrable. This was established by M. Gaudin in 1972gaudin and later by B. Jurcojurco who used a different approach. First explicit nonlinear solutions for the mean-field dynamics of the model (1,3) were constructed in Refs. andreev, ; barankov, . Interestingly, the initial conditions for which these solutions occur are such that the dynamics of the model (1,3) reduces to that of the Dicke model (4), i.e. it can be described in terms of a single collective spin coupled to the bosonic field. Similar solutions for the mean-field dynamics of the BCS model were discovered by R. A. Barankov et. al. in Ref. barankov1, and in an unpublished work by V. S. Shumeiko shumeiko .
Here we solve for the dynamics of the model (1,3) in the mean-field approximation for arbitrary initial conditions. We also derive a full set of conservation laws for the mean-field dynamics. The typical evolution of the system is quasi-periodic with $`n+1`$ incommensurate frequencies. However, for certain special choices of initial conditions the dynamics is characterized by $`m+1<n+1`$ incommensurate frequencies and can be described in terms of $`m`$ collective spins coupled to the bosonic field.
Our approach to the mean-field dynamics of the fermion-oscillator model (1) employs the method of separation of variables. This method was suggested by I. V. Komarovkomarov and later developed by E. Sklyanin and V. Kuznetsovsklyanin ; vadim . It allows us to derive a full set of integrals of motion for both the quantum model (1,3) and its classical counterpart. We also use it to derive and analyze equations of motion in terms of separation variables.
Upon the replacement of the quantum bosonic field $`\widehat{b}`$ with its time-dependent expectation value, the Hamiltonian (1) becomes similar to the mean-field BCS Hamiltonian. In this analogy, $`g\widehat{b}(t)`$ plays the role of the BCS gap $`\mathrm{\Delta }(t)`$. The difference is that in the case of the BCS model $`\mathrm{\Delta }(t)`$ is not an independent dynamical variable, but is related to $`\widehat{c}_j(t)\widehat{c}_j(t)=\widehat{K}_j^{}(t)`$ by the self-consistency condition $`\mathrm{\Delta }(t)=g_j\widehat{K}_j^{}(t)`$. Nevertheless, it turns out to be possible to map the mean-field dynamics of the fermion-oscillator model (1) with $`n`$ energy levels onto the corresponding BCS problem with $`n+1`$ levels. This mapping is facilitated by the variable separation and enables us to obtain the general solution for the time dependence of the expectation values $`\widehat{𝐊}_j(t)`$ and $`\widehat{b}(t)`$ in terms of hyperelliptic functionstheta ; mumford using the known solution for the mean-field dynamics of the BCS modelBCS ; shortBCS .
The rest of the paper is organized as follows. In section II, we derive the classical Hamiltonian that governs the mean-field dynamics of the quantum model (1,3). We show that the dynamics of the resulting classical model is integrable and derive a full set of its conservation laws in section III. In section IV, we perform a transformation to a new set of variables, which facilitates the solution of the equations of motion. The mapping to the corresponding BCS problem and the general solution for the mean-field dynamics of the model (1,3) are derived in section VI. Finally, section VII is devoted to an interesting class of particular solutions that include mean-field equilibrium states and special solutions of Refs. andreev ; barankov ; dicke ; BP .
## II Classical Model
Here we derive the classical Hamiltonian model that describes the mean-field dynamics of the quantum model (1,3).
We start with the Heisenberg equations of motion for the spin-oscillator model (3).
$$\begin{array}{c}\dot{\widehat{𝐊}}_j=i[\widehat{H},\widehat{𝐊}_j]=\widehat{𝐁}_j\times \widehat{𝐊}_j\widehat{𝐁}_j=(2g\widehat{b}_x,2g\widehat{b}_y,2ϵ_j),\hfill \\ \\ \dot{\widehat{b}}=i[\widehat{H},\widehat{b}]=i\omega \widehat{b}ig\widehat{T}_{}\widehat{𝐓}=\underset{q=0}{\overset{n1}{}}\widehat{𝐊}_q,\hfill \end{array}$$
(5)
where the operators $`\widehat{b}_x`$ and $`\widehat{b}_y`$ are defined by
$$\widehat{b}=\widehat{b}_xi\widehat{b}_y\widehat{b}^{}=\widehat{b}_x+i\widehat{b}_y.$$
In the regime when the bosonic mode is macroscopically populated, we can replace operators $`\widehat{b}(t)`$ and $`\widehat{b}^{}(t)`$ in Eqs. (5) with $`c`$-numbers: $`b(t)=\widehat{b}(t)`$ and $`\overline{b}(t)=\widehat{b}^{}(t)`$, where $`\mathrm{}`$ stands for the time-dependent quantum-mechanical expectation value. After this replacement, Eqs. (5) become linear in operators. Taking their quantum-mechanical expectation value, we obtain for $`𝐬_j(t)=\widehat{𝐊}_j(t)`$
$$\begin{array}{c}\dot{𝐬}_j=𝐁_j\times 𝐬_j𝐁_j=(2gb_x,2gb_y,2ϵ_j),\hfill \\ \\ \dot{b}=i\omega bigJ_{}𝐉=\underset{j=0}{\overset{n1}{}}𝐬_j.\hfill \end{array}$$
(6)
An important observation is that Eqs. (6) are Hamiltonian equations for the following classical model:
$$H=\underset{j=0}{\overset{n1}{}}2ϵ_js_j^z+\omega \overline{b}b+g\underset{j=0}{\overset{n1}{}}\left(\overline{b}s_j^{}+bs_j^+\right),$$
(7)
where $`s_j^\pm =s_j^x\pm is_j^y`$. This Hamiltonian governs $`n`$ classical spins interacting with a classical harmonic oscillatorclass .
Classical spins are similar to the angular momentum and have the same Poisson brackets
$$\begin{array}{c}\{s_j^a,s_j^b\}=\epsilon _{abc}s_j^c,\hfill \\ \\ \{b,\overline{b}\}=i\{b_x,b_y\}=\frac{1}{2},\hfill \end{array}$$
(8)
where $`a`$, $`b`$, and $`c`$ stand for the spatial indexes $`x`$, $`y`$, and $`z`$. All other Poisson brackets between components of $`𝐬_j`$ and $`b`$ vanish. Equations (8) determine the fundamental Poisson brackets in our problem in a sense that Poisson brackets between any other pair of dynamical variables (functions of $`b,\overline{b}`$ and components of $`𝐬_j`$) can be obtained from Eqs. (8) using the standard properties of Poisson bracketslandau . For example, Eqs. (6) can be obtained from the Hamiltonian equations of motion $`\dot{𝐬}_j=\{H,𝐬_j\}`$ and $`\dot{b}=\{H,b\}`$ using Eqs. (8) and (7).
Equations of motion (6) conserve the length of each spin, $`𝐬_j^2=\text{const}`$. This fact can also be viewed as a property of the brackets (8) – since $`𝐬_j^2`$ Poisson-commutes with all other dynamical variables, its bracket with the Hamiltonian (7) also vanishes. In the mean-field approximation eigenstates of the Hamiltonian (3) are product states. In a pure state of a spin-1/2, there is always an axis, $`𝐧`$, such that the projection of the spin onto it is $`1/2`$, i.e. $`\widehat{𝐊}_j𝐧=1/2`$. In this case, $`𝐬_j^2=\widehat{𝐊}_j^2=1/4`$. Therefore, if the system was in a product state at $`t=0`$
$$𝐬_j^2=(s_j^z)^2+s_j^+s_j^{}=\frac{1}{4}.$$
(9)
If a number of orbitals $`ϵ_j`$ in (7) are degenerate, the magnitude of their total classical spin $`_{ϵ_j=ϵ}𝐬_j`$ is conserved by the equations of motion (6). In this case, one can replace the corresponding spins $`𝐬_j`$ with a single classical spin of a larger magnitude,
$$\underset{ϵ_j=ϵ}{}ϵ_j𝐬_j=ϵ𝐒_ϵ,\text{where}𝐒_ϵ=\underset{ϵ_j=ϵ}{}𝐬_j,$$
(10)
in the Hamiltonian (7) and sum over nondegenerate orbitals only. Below we will assume whenever necessary that such a replacement has been made, i.e. that orbitals $`ϵ_j`$ are nondegenerate.
Finally, comparing Hamiltonians (3) and (7) and brackets (8) to the corresponding quantum commutators, we note that the mean-field approximation is equivalent to the standard procedure of going from quantum to classical mechanics by replacing commutators with Poisson brackets, $`i[A,B]\{A,B\}`$ (here $`\mathrm{}=1`$).
## III Integrability
Here we demonstrate that the classical model (7), as well as its quantum counterparts (3) and (1), are integrable and introduce a useful tool to analyze and solve for their dynamics.
The Hamiltonian (7) depends on $`2(n+1)`$ dynamical variables: 2 angles for each of $`n`$ spins plus the coordinate and the momentum of the oscillatorclass . Therefore, its phase space is $`2(n+1)`$-dimensional and the number of the degrees of freedom (the number of generalized coordinates) is $`n+1`$. To show that the classical model (7) is integrable, we have to show that it has $`n+1`$ independent integrals of motion (see e.g. Refs. arnold, ; landau, ).
Consider the following vector-functions of an auxiliary parameter $`u`$:
$$\begin{array}{c}𝐋_j(u)=\frac{𝐬_j}{uϵ_j}j=0,\mathrm{},n1\hfill \\ \\ 𝐋_n(u)=\frac{1}{g^2}\left(\begin{array}{c}2gb_x\\ 2gb_y\\ 2u\omega \end{array}\right).\hfill \end{array}$$
It follows from the above definitions and Eqs. (8) that components of $`𝐋_k(u)`$ have the following Poisson brackets:
$$\{L_k^a(v),L_k^b(w)\}=\frac{\epsilon _{abc}}{vw}\left(L_k^c(v)L_k^c(w)\right)k=0,\mathrm{},n$$
(11)
Because components of different $`𝐋_j(u)`$ Poisson-commute with each other, their sum
$$𝐋(u)=\frac{1}{g^2}\left(\begin{array}{c}2gb_x\\ 2gb_y\\ 2u\omega \end{array}\right)+\underset{j=0}{\overset{n1}{}}\frac{𝐬_j}{uϵ_j}$$
(12)
also satisfies Eqs. (11). One can check using only Eq. (11) that the square of this vector commutes with itself at any two values, $`v`$ and $`w`$, of the auxiliary parameter, i.e.
$$\{𝐋^2(v),𝐋^2(w)\}=0.$$
(13)
The function $`𝐋^2(u)`$ acts as a generating function for the model (7) and its integrals of motion. Indeed, evaluating $`𝐋^2(u)`$ from Eq. (12), we obtain
$$𝐋^2(u)=\frac{(2u\omega )^2}{g^4}+\frac{4H_n}{\omega g^2}+\underset{j=0}{\overset{n1}{}}\left[\frac{2H_j}{g^2(uϵ_j)}+\frac{𝐬_j^2}{(uϵ_j)^2}\right].$$
(14)
where
$$H_j=(2ϵ_j\omega )s_j^z+g(\overline{b}s_j^{}+bs_j^+)+g^2\underset{kj}{}\frac{𝐬_j𝐬_k}{ϵ_jϵ_k},$$
(15)
$$H_n=\omega \left(\overline{b}b+J_z\right)𝐉=\underset{j=0}{\overset{n1}{}}𝐬_j.$$
(16)
Hamiltonians $`H_j`$ defined by Eqs. (15,16) Poisson-commute with each other
$$\{H_j,H_k\}=0\text{for all }j\text{ and }k.$$
This follows directly from the fact that Eq. (13) holds for any $`v`$ and $`w`$. On the other hand, it is straightforward to verify that the spin-oscillator Hamiltonian (7) is a linear combination of $`H_j`$
$$H=\underset{j=0}{\overset{n}{}}H_j.$$
(17)
Since $`H_j`$ Poisson-commute with each other, the Hamiltonian (7) commutes with all $`n+1`$ Hamiltonians $`H_j`$. Therefore, $`H_j`$ given by Eqs. (15,16) are integrals of motion for the classical spin-oscillator model, i.e. this model is integrable.
The above construction based on Eqs. (11) and (12) constitutes the so-called Lax vector representation for the spin-oscillator model (7). Its main advantages are that it provides a powerful tool to study the dynamics and can also be extended to quantum models. In particular, note that integrals (15) and (16) can be trivially generalized to the quantum case
$$\begin{array}{c}\widehat{H}_j=(2ϵ_j\omega )\widehat{K}_j^z+g(\widehat{b}^{}\widehat{K}_j^{}+\widehat{b}\widehat{K}_j^+)+g^2\underset{jk}{}\frac{\widehat{𝐊}_j\widehat{𝐊}_k}{ϵ_jϵ_k},\hfill \\ \\ \widehat{H}_n=\omega \left(\widehat{b}^{}\widehat{b}+\widehat{T}_z\right)\widehat{𝐓}=\underset{q=0}{\overset{n1}{}}\widehat{𝐊}_q.\hfill \end{array}$$
Operators $`\widehat{H}_j`$ pairwise commute with each other and the spin-oscillator Hamiltonian (3), which is their linear combination as in Eq. (17).
The same construction with classical dynamics in the mean-field approximation and Lax vector for resulting classical models applies to a number of other models. In particular, a closely related model is the BCS model
$$\widehat{H}_{BCS}=\underset{j,\sigma }{}ϵ_j\widehat{c}_{j\sigma }^{}\widehat{c}_{j\sigma }g\underset{j,q}{}\widehat{c}_j^{}\widehat{c}_j^{}\widehat{c}_q\widehat{c}_q.$$
(18)
Notations here are the same as in Eq. (1). In terms of Anderson’s spins (see above Eq. (3)), the Hamiltonian reads
$$\widehat{H}_{BCS}=\underset{j=0}{\overset{n}{}}2ϵ_j\widehat{K}_j^zg\underset{j,q}{}\widehat{K}_j^+\widehat{K}_q^{}.$$
(19)
The usual BCS mean-field is equivalent to the procedure (Eqs. (5) and (6)) we used to derive the classical Hamiltonian (7), only now the role of $`\widehat{b}`$ is played by $`_j\widehat{K}_j^{}`$ (see e.g. Refs. barankov1, ; BCS, ). The mean-field dynamics is described by the following classical Hamiltonian:
$$H_{BCS}=\underset{j=0}{\overset{n}{}}2ϵ_js_j^zg\underset{j,q}{}s_j^+s_q^{}$$
(20)
with the Lax vector
$$𝐋_{BCS}(u)=\frac{\widehat{𝐳}}{g}+\underset{j}{}\frac{𝐬_j}{uϵ_j},$$
(21)
where $`\widehat{𝐳}`$ is a unit vector along the $`z`$-axis. The Lax vector $`𝐋_{BCS}(u)`$ has the same properties (11) and (13) as before and generates integrals of motion for the BCS model.
## IV Separation of variables
In this section we perform a transformation to a new set of variables, which facilitates the solution of the equations of motion (6). The choice of new variables naturally follows from the Lax vector construction of the previous section. The new variables are canonical and also separate the Hamilton-Jacobi equations for the classical Hamiltonian (7) in the usual sensearnold ; landau .
We define $`n`$ variables $`u_j`$ as zeros of
$$L_{}(u)=L_x(u)iL_y(u)=\frac{2b}{g}+\underset{j=0}{\overset{n1}{}}\frac{s_j^{}}{uϵ_j}.$$
(22)
Note that the nominator of this expression is a polynomial of degree $`n`$ and therefore there are $`n`$ roots $`u_j`$ with $`j=0,\mathrm{},n1`$. The coefficients of this polynomial are functions of the dynamical variables $`s_j^{}`$ and $`b`$. Accordingly, its roots are also functions of $`s_j^{}`$ and $`b`$ and therefore define a new set of dynamical variables.
Variables $`u_j`$ play a role of canonical coordinates for the classical Hamiltonian (7). The corresponding momenta are defined as $`v_j=L_z(u_j)`$. Thus, we have
$$L_{}(u_j)=0,v_j=L_z(u_j),j=0,\mathrm{},n1.$$
(23)
Because our system has $`2(n+1)`$ degrees of freedom, we need two additional variables $`u_n`$ and $`v_n`$, which can be introduced as
$$u_n=b,v_n=\frac{H_n}{\omega b}=\overline{b}+\frac{J_z}{b},$$
(24)
where $`H_n`$ is given by Eq. (16).
Separation variables $`(u_j,v_j)`$ are canonical, i.e.
$$\{u_j,u_k\}=0\{v_j,v_k\}=0\{v_j,u_k\}=i\delta _{jk}j,k=0,\mathrm{},n.$$
(25)
The first relation in Eq. (25) follows from the fact that by Eqs. (23) and (24) variables $`u_j`$ depend only on mutually Poisson-commuting variables $`s_k^{}`$ and $`b`$. The second relation follows from $`\{L_z(v),L_z(w)\}=\{L_z(v),b\}=0`$. To derive Poisson brackets between $`v_j`$ and $`u_k`$ for $`j,k=0,\mathrm{},n1`$; we use the following equation obtained from Eq. (11)
$$\{L_z(v),L_{}(w)\}=\frac{i}{vw}\left(L_{}(v)L_{}(w)\right).$$
(26)
Evaluating this expression at $`v=u_j`$ and $`w=u_k`$, we obtain
$$\{L_z(u_j),L_{}(w)\}_{w=u_k}=\underset{wu_k}{lim}\frac{i}{u_jw}L_{}(w)=iL_{}^{}(u_k)\delta _{jk}.$$
(27)
where $`L^{}(u)=L/u`$. On the other hand,
$$\{L_z(u_j),L_{}(w)\}_{w=u_k}=\{L_z(u_j),L_{}(u_k)\}L_{}^{}(u_k)\{L_z(u_j),u_k\}=L_{}^{}(u_k)\{L_z(u_j),u_k\},$$
(28)
where we used $`L_{}(u_k)=0`$. Comparing Eqs. (27) and (28), we obtain the last relation in Eq. (25) for $`j,k=0,\mathrm{},n1`$.
The original dynamical variables $`𝐬_j=\widehat{𝐊}_j(t)`$ and $`b=\widehat{b}(t)`$ that we are interested in, can be explicitly expressed in terms of $`u_j`$ (the inverse map). Indeed, note that $`L_{}(u)`$ has its zeros at $`u=u_j`$ and poles at $`u=ϵ_j`$. Therefore, using Eqs. (23) and (22), we can write it as
$$L_{}(u)=\frac{2b}{g}+\underset{j=0}{\overset{n1}{}}\frac{s_j^{}}{uϵ_j}=\frac{2b}{g}\frac{_k(uu_k)}{_j(uϵ_j)},$$
(29)
where we have also used $`L_{}(u)=2b/g+O(1/u)`$ for large $`u`$. Equating residues at $`u=ϵ_j`$ and $`u=\mathrm{}`$, we obtain
$$s_j^{}=\frac{2b}{g}\frac{_k(ϵ_ju_k)}{_{lj}(ϵ_jϵ_l)},$$
(30)
$$J_{}=\underset{j}{}s_j^{}=\frac{2b}{g}\underset{j=0}{\overset{n1}{}}(ϵ_ju_j).$$
(31)
Similarly, using
$$L_z(u_j)=v_jL_z(u)=\frac{2u\omega }{g^2}+O(1/u),$$
one can express $`L_z(u)`$ in terms of $`u_j`$. The $`z`$-components of classical spins, $`s_j^z`$, are residues of $`L_z(u)`$ at $`u=ϵ_j`$ (see Eq. (12))
$$s_j^z=s_j^{}\left[\frac{ϵ_j\omega /2+_k(u_kϵ_k)}{bg}+\underset{k}{}\frac{v_k}{(ϵ_ju_k)L_{}^{}(u_k)}\right].$$
(32)
It also follows from Eq. (23) that $`v_j^2=𝐋^2(u_j)`$. This allows us to express variables $`v_j`$ through $`u_j`$ and the integrals of motion $`H_j`$
$$v_j^2=\frac{(2u_j\omega )^2}{g^4}+\frac{4H_n}{\omega g^2}+\underset{k=0}{\overset{n1}{}}\left[\frac{2H_k}{g^2(u_jϵ_k)}\frac{𝐬_k^2}{(u_jϵ_k)^2}\right]j=0,\mathrm{},n1$$
(33)
$$v_n=\frac{H_n}{\omega u_n}=\frac{H_n}{\omega b}.$$
(34)
Thus, to determine the evolution of $`𝐬_j(t)=\widehat{𝐊}_j(t)`$ and $`b(t)=\widehat{b}(t)`$ we only need to derive and solve the equations of motion for $`u_j(t)`$.
## V Equations of motion for separation variables
In order to derive the equations of motion for separation variables $`u_j`$, we first determine the brackets $`u_{l,k}=\{H_k,u_l\}`$ for $`l,k=0,\mathrm{},n1`$ using Eqs. (33) and (25). This is done by taking Poisson brackets of both sides of Eq. (33) with $`u_l`$ and solving the resulting system of linear equations for $`u_{l,k}`$.
As soon as $`u_{l,k}`$ are found in terms of $`u_j`$ and $`v_j`$, we can use the expression (17) for the Hamiltonian in terms of $`H_j`$ to determine $`\dot{u}_j`$
$$\dot{u}_j=\{H,u_j\}=\underset{k}{}u_{j,k}.$$
(35)
Taking Poisson brackets of both sides of Eqs. (33) and (34) with $`u_l`$, we obtain
$$ig^2v_j\delta _{jl}=\underset{k=0}{\overset{n1}{}}\frac{u_{l,k}}{u_jϵ_k}j,l=0,\mathrm{}n1.$$
(36)
In order to determine $`u_{n,k}`$ and $`u_{l,k}`$ for $`l,k=0,\mathrm{},n1`$, we need to invert the matrix
$$A_{jk}=\frac{1}{u_jϵ_k},$$
(37)
which is called the Cauchy matrix.
This can be done with the help of the following identity coming from partial fraction decomposition of the LHS:
$$\frac{_{l=0}^{n1}(uu_l)}{(uu_p)_{l=0}^{n1}(uϵ_l)}=\underset{j=0}{\overset{n1}{}}\frac{1}{uϵ_j}\frac{_{lp}(ϵ_ju_l)}{_{lj}(ϵ_jϵ_l)}.$$
(38)
To verify this identity note that both sides have the same poles and that the residues at these poles coincide. Substituting $`u=u_q`$ in Eq. (38), one derives
$$\delta _{qp}=\underset{j=0}{\overset{n1}{}}\frac{1}{u_qϵ_j}\frac{_{l=1}^n(u_pϵ_l)}{_{lj}(ϵ_jϵ_l)}\underset{lp}{}\frac{ϵ_ju_l}{u_pu_l}.$$
(39)
Hence
$$\left(A^1\right)_{jk}=\frac{_{l=1}^n(u_kϵ_l)}{_{lj}(ϵ_jϵ_l)}\underset{lk}{}\frac{ϵ_ju_l}{u_ku_l}.$$
(40)
It follows from Eqs. (36) and (30) that
$$u_{j,k}=ig^2v_j\frac{_{l=0}^{n1}(u_jϵ_l)}{_{lj}(u_ju_l)}\frac{gs_k^{}}{2b(ϵ_ku_j)}.$$
(41)
Finally, noticing that $`0=L_{}(u_j)=2b/g+_ks_k^{}/(u_jϵ_k)`$ and using Eqs. (35) and (33), we derive
$$\begin{array}{c}\dot{u}_j=\frac{2i\sqrt{Q_{2n+2}(u_j)}}{_{mj}(u_ju_m)}j=0,\mathrm{},n1\hfill \\ \\ \dot{b}=2ib\left(\frac{\omega }{2}+\underset{j=0}{\overset{n1}{}}ϵ_j\underset{k=0}{\overset{n1}{}}u_k\right),\hfill \end{array}$$
(42)
where the spectral polynomial $`Q_{2n+2}(u)`$ is defined as
$$Q_{2n+2}(u)=g^4𝐋^2(u)\underset{j=0}{\overset{n1}{}}(uϵ_j)^2.$$
(43)
By Eq. (14), the coefficients of $`Q_{2n+2}(u)`$ depend only on the integrals of motion $`H_j`$ and parameters $`ϵ_j`$ and $`g`$. The equation of motion for $`u_n=b`$ is obtained by substituting Eq. (31) into Eq. (6).
An almost identical derivation of equations of motion for separation variables can be performed for the mean-field BCS model (20) using its vector Lax representation (21) leading to the following equations of motion:
$$\begin{array}{c}\dot{u}_j=\frac{2i\sqrt{Q_{2n+2}^{BCS}(u_j)}}{_{mj}(u_ju_m)}j=0,\mathrm{},n1\hfill \\ \\ \dot{J}_{}=2iJ_{}\left(gJ_z+\underset{j=0}{\overset{n}{}}ϵ_j\underset{k=0}{\overset{n1}{}}u_k\right),\hfill \end{array}$$
(44)
where the spectral polynomial for the BCS model is defined as
$$Q_{2n+2}^{BCS}(u)=\frac{1}{g^2}𝐋_{BCS}^2(u)\underset{j=0}{\overset{n}{}}(uϵ_j)^2.$$
(45)
## VI The solution
Here we obtain the general solution for the mean-field dynamics of the fermion-oscillator (spin-oscillator) model (1,3) by mapping it onto the corresponding BCS model. The solution therefore can be read off from the known general solution of the mean-field BCS problemBCS .
By comparing equations of motion (42) for the spin-oscillator model with those for BCS model (44), we observe that they coincide upon a replacement
$$\underset{j=0}{\overset{n}{}}ϵ_j+gJ_z\underset{j=0}{\overset{n1}{}}ϵ_j+\frac{\omega }{2},J_{}b,Q_{2n+2}^{BCS}(u)Q_{2n+2}(u).$$
(46)
Thus, the mean-field dynamics of the spin-oscillator model (3) with $`n`$ spins and a bosonic field coincides with that of the BCS model (19) with $`n+1`$ spins. This allows us to immediately write down the solution for the time-dependent averages $`𝐬_j(t)=\widehat{𝐊}_j(t)`$ and $`b(t)=\widehat{b}(t)`$
$$b(t)=\left[J_{}\right]_{BCS},[𝐬_j(t)]_{\text{Dicke}}=\left[𝐬_j(t)\right]_{BCS}.$$
(47)
The explicit form of $`\left[J_{}\right]_{BCS}`$ and $`[𝐬_j(t)]_{BCS}`$ in terms of hyperelliptic functions can be found in Ref. BCS, (see Eqs. (3.22 – 3.24) of Ref. BCS, ).
The dynamics of $`𝐬_j(t)=\widehat{𝐊}_j(t)`$ and $`b(t)=\widehat{b}(t)`$ is typical of a classical integrable system with $`n+1`$ degrees of freedomarnold (recall that $`n`$ is the number of non-degenerate single particle levels $`ϵ_j`$ in the fermion-oscillator model (1)). The typical motion is quasi-periodic with $`n+1`$ incommensurate periods. For example, a Fourier transform of $`b(t)=\widehat{b}(t)`$ shows $`n+1`$ basic frequencies. The system uniformly (ergodically) explores the invariant torus – the $`n+1`$-dimensional portion of the $`2n+2`$-dimensional phase space allowed by the integrals of motion (15) and (16), returning arbitrarily close to the initial point at irregular intervals.
We found that in the thermodynamic limit $`n\mathrm{}`$ the solution simplifies for most physical initial conditions. The return time diverges in this limit. The dynamics of $`b(t)`$ is particularly simple – it decays, typically as a power-law, to a steady state where $`|b(t)|`$ is either constant or is characterized by a few independent frequenciesunpub . However, the motion of the spin system still contains a continuum of frequencies. The final steady state of $`|b(t)|`$ depends on the initial conditions.
## VII Few spin solutions
The evolution described in the previous section occurs for most initial conditions and is stable against small perturbations of the Hamiltonian by the KAM theorem even if these perturbations destroy the integrability. However, there always exists a set of points of measure zero in the phase space, where the motion is characterized by only a few incommensurate frequencies, while the stability is not guaranteed arnold . Below we consider a family of such solutions, which we call $`m`$-spin solutions with $`m<n`$. The reason is that in these cases the dynamics of $`n`$ spins and the bosonic field in Hamiltonian (3) degenerates to that of $`m<n`$ spins and the bosonic field.
Few spin solutions are constructed by choosing integrals of motion $`H_j`$ (i.e. the initial conditions) so that $`2(nm)`$ roots of the spectral polynomial $`Q_{2n+2}(u)`$ defined in Eq. (14) become double degenerate. Suppose $`u=E_0`$ is a double rootreal of $`Q_{2n+2}(u)`$ and note that, since $`Q_{2n+2}(E_0)=𝐋^2(E_0)=0`$, setting e.g. $`u_{n1}(t)=E_0`$ solves the equation of motion (42) for the variable $`u_{n1}`$. Therefore, we can ”freeze” one of the separation variables in this root. Then $`L_{}(E_0)=0`$, and it follows from $`L_z^2(u)+L_{}(u)L_+(u)=𝐋^2(u)`$ that $`L_z(E_0)=0`$. Thus, one can factor out $`(uE_0)`$ from all components of the Lax vector and show that it is proportional to the Lax vector of the spin-oscillator model with $`n1`$ spins. Similarly, if there are 2 pairs of double degenerate roots, the number of spins reduces to $`n2`$ etc. This procedure is followed in detail in Ref. BCS, using a different method.
First, let us consider the general case when the number of spins effectively reduces from $`n`$ to $`m<n`$. We have
$$𝐋(u)=\left(1+\underset{j}{}\frac{d_j}{uϵ_j}\right)𝐋_𝐭(u),$$
(48)
where $`d_j`$ are time-independent constants and $`𝐋_𝐭(u)`$ is the Lax vector of the $`m`$-spin problem
$$𝐋_𝐭(u)=\frac{1}{g^2}\left(\begin{array}{c}2gb_x\\ 2gb_y\\ 2u\omega ^{}\end{array}\right)+\underset{k=0}{\overset{m1}{}}\frac{𝐭_k}{u\mu _k}.$$
(49)
The $`m`$-spin system has its own $`m`$ arbitrary ”energy levels” $`\mu _k`$. Its dynamics is governed by the spin-oscillator Hamiltonian (7) for $`m`$ spins with new parameters replacing $`ϵ_j`$ and $`\omega `$.
$$H_m=\underset{k=0}{\overset{m1}{}}2\mu _kt_k^z+\omega ^{}\overline{b}b+g\underset{k=0}{\overset{m1}{}}\left(\overline{b}t_k^{}+bt_k^+\right).$$
(50)
Matching the residues at poles at $`u=ϵ_j`$ on both sides of Eq. (48), we express the original spins $`𝐬_j`$ in terms of collective spins $`𝐭_k`$
$$𝐬_j=d_j𝐋_𝐭(ϵ_j).$$
(51)
Constants $`d_j`$ are determined from the above equation using $`𝐬_j^2=1/4`$. We have
$$d_j=\frac{|𝐬_j|e_j}{\sqrt{𝐋_𝐭^2(ϵ_j)}}=\frac{e_j}{2\sqrt{𝐋_𝐭^2(ϵ_j)}}e_j=\pm 1.$$
(52)
Finally we have to match the residues at $`u=\mu _k`$ and the $`u\mathrm{}`$ asymptotic. This leads to the following $`m+1`$ equations
$$\begin{array}{c}1+\underset{j=0}{\overset{n1}{}}\frac{d_j}{\mu _kϵ_j}=0k=0,\mathrm{},m1\hfill \\ \\ \omega =\omega ^{}2\underset{j=0}{\overset{n1}{}}d_j.\hfill \end{array}$$
(53)
Equations (53) constrain the lengths of new spins $`𝐭_k`$, i.e. the coefficients of the spectral polynomial $`Q_{2m+2}(u)=g^4𝐋_𝐭^2(u)_{k=0}^{m1}(u\mu _k)^2`$ of the $`m`$-spin system (see Eq. (14)). Indeed, using Eq. (52), one can cast the constraints Eq (53) into the form
$$\underset{j=0}{\overset{n1}{}}\frac{e_jϵ_j^{l1}}{\sqrt{Q_{2m+2}(ϵ_j)}}=\frac{2}{g^2}\delta _{lm}l=1,\mathrm{},m$$
(54)
$$\omega ^{}=\omega +2\underset{k=0}{\overset{m1}{}}\mu _k\underset{j=0}{\overset{n1}{}}\frac{g^2e_jϵ_j^m}{\sqrt{Q_{2m+2}(ϵ_j)}}.$$
(55)
To obtain $`m`$-spin solutions explicitly, one has to choose parameters $`\mu _k`$, resolve the $`m+1`$ constraints (54) and (55) for the lengths of spins $`𝐭_k`$ and frequency $`\omega ^{}`$, and solve for the dynamics of the $`m`$-spin Hamiltonian (50). The dynamics can be obtained from the general solution (47) by replacing $`nm`$ and the set of $`\{ϵ_j\}`$ with $`\{\mu _k\}`$.
Let us illustrate the construction of $`m`$-spin solutions by considering the cases $`m=0`$ and $`m=1`$ in more detail.
m=0. The 0-spin solutions correspond to the mean-field eigenstates of quantum Hamiltonians (1) and (3). They are equilibrium states for the classical Hamiltonian (7).
For $`m=0`$ the dynamics is governed by the Hamiltonian $`H_0=\omega ^{}\overline{b}b`$ with an obvious solution $`b=b_0e^{i(\omega ^{}t+\varphi )}`$, where $`b_0=|b(t)|=\text{const}`$. Further, using Eqs. (49), (51), and (52), we obtain
$$𝐋_𝐭^2(u)=\frac{1}{g^4}\left(4g^2b_0^2+(2u\omega ^{})^2\right),$$
(56)
$$\widehat{c}_j\widehat{c}_j=s_j^{}=\frac{e_jgb}{\sqrt{4g^2b_0^2+(2ϵ_j\omega ^{})^2}},\widehat{n}_j1=2s_j^z=\frac{2e_j(2ϵ_j\omega ^{})}{\sqrt{4g^2b_0^2+(2ϵ_j\omega ^{})^2}}.$$
(57)
The frequency $`\omega ^{}`$ and the gap $`\mathrm{\Delta }_0=gb_0`$ have to satisfy a single constraint (55) that now reads
$$\omega ^{}=\omega \underset{j=0}{\overset{n1}{}}\frac{g^2e_j}{\sqrt{4g^2b_0^2+(2ϵ_j\omega ^{})^2}}.$$
(58)
We see from Eq. (57), that all spins rotate uniformly around the $`z`$-axis with a frequency $`\omega ^{}`$. This rotation can be eliminated, i.e. $`\omega ^{}`$ can be set to zero, by an appropriate choice of the chemical potential. Then, Eq. (58) is the analog of the BCS gap equation. The configuration of spins (57) corresponds to the product BCS wavefunction. The latter can be straightforwardly reconstructed from the knowledge of $`𝐬_j=\widehat{𝐊}_j`$ as fixing the average of spin-1/2 uniquely fixes its quantum state. The choice of signs $`e_j=+1`$ for all $`j`$ corresponds to the ground state, while choosing one of the signs to be $`1`$ corresponds to an excited state – pair excitation of the fermionic condensate.
There is also an important class of 0-spin solutions that is obtained by setting $`b_0=0`$. We see from Eq. (57) that in this case all spins are along the $`z`$-axis, $`s_j^z=e_j/2=\pm 1/2`$. The constraint (58) is now irrelevant, because all $`xy`$-components vanish. The choice of signs $`e_j=\text{sign }ϵ_j`$, where $`ϵ_j`$ are counted from the Fermi level, yields the Fermi ground state. Other choices correspond to the excitations of the Fermi gas. All these states are stationary with respect to the mean-field dynamics. For a finite number of pairs, they are non-stationary with respect to the quantum Hamiltonian (1). Thus, their initial evolution is entirely governed by quantum corrections (cf. Refs. BP, ; leggett, ; shortBCS, ).
m=1. In this case the dynamics reduces to that of a single collective spin $`𝐭`$ coupled to an oscillator, i.e. it is governed by the classical counterpart of the Dicke model (4)
$$H_1=2\mu t_z+\omega ^{}\overline{b}b+g(\overline{b}t_{}+bt_+).$$
(59)
Using Eq. (57), one can express original spins in terms of the collective spin and the bosonic field
$$s_j^{}=\frac{e_j(ϵ_j\mu )b+ge_jt_{}}{2g\sqrt{Q_4(ϵ_j)}}s_j^z=\frac{e_j(2ϵ_j\omega ^{})(ϵ_j\mu )+g^2t_z}{2g^2\sqrt{Q_4(ϵ_j)}}.$$
(60)
To complete the construction of 1-spin solutions, we need to choose a positively defined polynomial $`Q_4(u)`$ so that it satisfies two constraints (54) (55) and solve for the dynamics of the classical Dicke model (59). In this case there is only one non-stationary separation variable $`u_0`$ and equations of motion (42) take the following form:
$$\begin{array}{c}\dot{u}_0^2+4Q_4(u_0)=0,\hfill \\ \dot{b}=ib(\omega +2\mu 2u_0).\hfill \end{array}$$
(61)
These equations can be solved in terms of elliptic functions. This is not surprising, since the model (59) is, in fact, equivalent to a spherical pendulum and its solution has already been obtained in Refs. dicke ; BP .
The 1-spin solutions were originally obtained in Refs. andreev ; barankov where they were used to describe the evolution beginning from a state infinitesimally close to the normal ground state.
## VIII Conclusion
In this paper we solved for the mean-field dynamics of the fermion-oscillator model (1). In the mean-field approximation, the problem reduces to a classical Hamiltonian model. We derived integrals of motion for both the classical (7) and the quantum (1) models and showed that the dynamics of the fermion-oscillator model (1) maps onto that of the BCS model. This was used to derive an explicit general solution for the mean-field dynamics of the fermion-oscillator model with an arbitrary finite number $`n`$ of degrees of freedom.
The typical dynamics is quasi-periodic with $`n`$ incommensurate frequencies. The system ergodically explores the part of the phase-space allowed by integrals of motion, returning arbitrarily close to its initial state at irregular time intervals. We have also constructed a class of particular, few-spin, solutions, for which the dynamics reduces to that of $`m`$ collective spins governed by the same classical model. The case $`m=0`$ corresponds to the mean-field eigenstates of the fermion-oscillator Hamiltonian (1).
An interesting problem is to obtain and describe the thermodynamic limit, $`n\mathrm{}`$, of the solution. For most physical initial conditions the dynamics considerably simplifies in this limitunpub . In particular, we found that $`\widehat{b}(t)`$ decays, typically as a power-law, to a steady state where $`|\widehat{b}(t)|`$ is either a constant (cf. Refs. galaiko, ; volkov, ; simons, ; leggett, ) or oscillates with few independent frequencies (cf. Fig. 2 in Refs. galperin, ; barankov1, ). The motion of the spin system still contains a continuum of frequencies. The steady state of $`|\widehat{b}(t)|`$ depends on the initial conditions.
It is also interesting to study quantum and finite temperature effects. Quantum corrections to the mean-field dynamics can become important when the bosonic mode is weakly populated BP . For example, normal states where $`\widehat{b}^{}\widehat{b}=\widehat{c}_j\widehat{c}_j=0`$ are stationary in the mean-field approximation (6), but are not stationary with respect to quantum dynamics (5). |
warning/0506/gr-qc0506140.html | ar5iv | text | # GAUGE APPROACH TO GRAVITATION AND REGULAR BIG BANG THEORY
## I I. Introduction
The Hot Bing Bang scenario built on a base of Friedmannian homogeneous isotropic cosmological models (HICM) of general relativity theory (GR) is the foundation of relativistic cosmology. Well known physical laws allow to describe physical processes in evolving Universe beginning from the time $`t_110^4s`$ from the Bing Bang corresponding to the start of cosmological expansion (in the case of Friedmannian HICM the moment $`t=0`$ corresponds to cosmological singularity), when the energy density of gravitating matter was comparable with nuclear density and the temperature $`T10^{11}K`$. Observable confirmation of theoretical predictions of primordial nucleosynthesis of light elements, explanation of large scale structure formation of the Universe are important achievements of the Hot Bing Bang theory. To describe gravitating matter and physical processes in evolving Universe at $`t<t_1`$ one has to use modern theory of fundamental physical interactions, unified models of particle physics, that is the object of cosmology of early Universe. The creation of inflationary paradigm, which permits to solve a number of problems of standard cosmological scenario, in particular, to explain the homogeneity and isotropy of the Universe at initial stages of cosmological expansion is important achievement of early Universe cosmology mc27 . The further progress in cosmology of early Universe is connected with further perfection of particle physics theory and also theory of gravitational interaction.
As it is known, GR has problems by description of gravitating matter at extreme conditions with extremely high energy densities, pressures, temperatures (problem of gravitational singularities). The Friedmannian HICM possess singular state in the past, in which some important physical characteristics of gravitating matter (energy density, temperature etc) diverge and the scale factor $`R(t)`$ of Robertson-Walker metrics vanishes. As a result the existence of Friedmannian HICM is limited in the time in the past. Corresponding problem - the problem of cosmological singularity (PCS) – is one of the most principal cosmological problems, which does not have general solution still. From physical point of view, the PCS is connected with the fact that gravitational interaction in the case of usual gravitating systems with positive values of energy density $`\rho `$ and pressure $`p`$ in the frame of GR as well as Newton s gravitation theory has the character of attraction, but not repulsion. Although in the case of gravitating systems with negative pressure (for example, massive and nonlinear scalar fields) the gravitational interaction in GR can have the repulsion character, however, the PCS can not be solved by taking such systems into account in the frame of GR mc4 ; mc5 . There were many attempts to resolve PCS in the frame of GR as well as other classical theories of gravitation (see mc1 ; mc15 and Refs given herein). A number of regular cosmological solutions was obtained in the frame of metric theories of gravitation and also other theories, in the frame of which gravitation is described by using more general geometry than the Riemannian one. In connection with this, note that the resolution of PCS means not only obtaining regular cosmological solutions, but also excluding singular solutions of cosmological equations, as a result generic feature of cosmological solutions has to be its regular character. Moreover, gravitation theory and cosmological equations have to satisfy the correspondence principle with GR in the case of usual gravitating systems with sufficiently small energy densities and weak gravitational fields excluding nonphysical solutions. The greatest part of existent attempts to resolve the PCS does not satisfy indicated conditions mc15 .
Usually the Universe evolution in the frame of standard cosmological scenario is considered beginning from the time $`t`$ greater than the Planckian time $`t_p10^{43}s`$,when the energy density $`\rho `$ and temperature $`T`$ were smaller than the Planckian ones: $`\rho _p910^{112}\frac{\text{J}}{\text{m}^3}`$ and $`T_p10^{32}K`$.Then according to general opinion quantum gravitational effects are negligible and gravitational field can be described classically. On the contrary, at Planckian conditions ($`\rho \rho _p`$, $`TT_p`$) one means that quantum gravitational effects must be essential, and possibly the PCS can be solved by taking these effects into account. So, in order to solve the PCS the idea of quantum birth of the Universe was introduced tryon ; fomin . There are several realizations of this idea in the frame of quantum cosmology by using Wheeler-DeWitt equation (see review mc11 ). The appearing closed micro-universe is transformed then into macro-universe by virtue of inflation. Although singular state with divergent energy density is absent in such model, the problem of the beginning in the time of the Universe remains. The theory of quantum micro-universe was developed further in the frame of loop quantum gravity. The space quantization permits to continue solution for micro-universe of closed and flat type to the past and obtain bouncing solution containing the compression stage before the expansion boj . Note, this scenario has some vaguenesses, in particular: if at compression stage the Universe was a macro-object filled by gravitating matter what physical factors lead to transformation of macro-universe into micro-universe before a bounce, what is a flat model with finite volume ?
The PCS was discussed also in the frame of other candidate for quantum gravity theory - string theory/ M-theory lid ; ven1 ; khory ; mc12 ; mc13 ; ven2 . Some cosmological scenarios with evolution stage before the Big Bang - Pre-Big Bang scenario ven1 and ekpyrotic scenario khory \- were proposed. Note that HICM in a four-dimensional Einsteinian frame, generally speaking, are singular lid . The obtaining of regular bouncing solutions is connected with violation of condition of energy density positivity for gravitating matter . So, in Ref. ven1 non-local negative scalar field potential is introduced, in Ref. ven2 a component with negative energy density is used. To build oscillating cosmological model, the specific negative scalar field potential was introduced in Refs. mc12 ; mc13 . Although the energy density does not diverge by transition from compression to expansion, the singularity connected with vanishing of the scale factor $`R(t)`$ remains.
As it was shown in a number of our papers (see mc15 and Refs given herein), the applying of gauge approach to gravitation permits to build regular Big Bang theory by satisfying positivity property of energy density for gravitating matter. The present paper is devoted to consideration of principal features of developing approach and to further analysis of regular inflationary cosmology. In Section 2 generalized cosmological Friedmann equations for HICM in the frame of gauge approach to gravitation are introduced. In Section 3 these equations are applied to analyze inflationary cosmological models filled by scalar fields and usual gravitating matter. As illustration of regular Big Bang theory in Section 4 particular inflationary cosmological model is discussed.
## II II. GAUGE APPROACH TO GRAVITATION AND GENERALIZED COSMOLOGICAL FRIEDMANN EQUATIONS
As it is known, the local gauge invariance principle is the basis of modern theory of fundamental physical interactions. The theory of electro-week interaction, quantum chromodynamics, Grand Unified models of particle physics were built by using this principle. From physical point of view, the local gauge invariance principle establishes the correspondence between certain important conserving physical quantities, connected according to the Noether’s theorem with some symmetries groups, and fundamental physical fields, which have as a source corresponding physical quantities and play the role of carriers of fundamental physical interactions. The applying of this principle to gravitational interaction leads, generally speaking, to generalization of Einsteinian theory of gravitation. Metric theories of gravitation including GR, in the frame of which the energy-momentum tensor is a source of gravitational field, can be introduced in the frame of gauge approach by using the localization of 4-parametric translation group mc16 <sup>1</sup><sup>1</sup>1Because in the frame of gauge approach the gravitational interaction is connected with space-time transformations, the gauge treatment to gravitation has essential differences in comparison with Yang-Mills fields connected with internal symmetries groups. As a result, there are different gauge treatments of gravitational interaction not discussed in this paper.. Because the localized translation group is, in fact, the group of general coordinate transformations, from this point of view the general covariance of GR plays dynamical role. At the same time the Lorentz group (group of tetrad Lorentz transformations) in GR and other metric theories of gravitation does not play any dynamical role from the point of view of gauge approach, because corresponding Noether s invariant in these theories is identically equal to zero mc17 . If one means that the Lorentz group plays the dynamical role in gauge field theory and the Lorentz gauge field exists in the nature, in this case we obtain with necessity the gravitation theory in the Riemann-Cartan space-time (see, for example, hehl ; hay ). Corresponding theory is known as Poincare gauge theory of gravitation (PGTG). Gravitational field variables in PGTG are the tetrad (translational gauge field) and Lorentz connection (Lorentz gauge field); corresponding field strengths are torsion and curvature tensors. As sources of gravitational field in PGTG are covariant generalizations of energy-momentum and spin tensors. Unlike gauge Yang-Mills fields, for which the Lagrangian is quadratic in the gauge field strengths, gravitational Lagrangian of PGTG can include also linear in curvature term (scalar curvature), which is necessary to satisfy the correspondence principle with GR.
The first attempt to apply the simplest PGTG - Einstein-Cartan theory - in order to solve the PCS was made in Refs. kop ; traut . By using some classical model for spinning matter, non-singular cosmological solutions were obtained. However, it was shown later, these solutions have model character and critically depend on spinning matter description; by another spin description (for example, by means of Dirac field) the cosmological singularity does not disappear. Moreover, because in the case of spinless matter the Einstein-Cartan theory is identical to GR, all singular solutions for Friedmannian HICM are exact solutions of Einstein-Cartan theory of gravitation. The next step to apply the PGTG in order to solve the PCS was made in Ref.mc14 . In the frame of PGTG based on general expression of gravitational Lagrangian $`L_G`$ including both a scalar curvature and different invariants quadratic in the curvature and torsion tensors gravitational equations for HICM were deduced, these equations lead to the following generalized cosmological Friedmann equations (GCFE) <sup>2</sup><sup>2</sup>2The structure of gravitational equations of PGTG for HICM does not depend on detailed form of quadratic part of gravitational Lagrangian $`L_G`$. As a result, these equations contain two independent parameters $`f`$ and $`a`$, which are linear combinations of coefficients at terms of $`L_G`$ quadratic in the curvature and torsion tensors. The mathematical requirement for cosmological equations for HICM to be differential equations of the same order as in GR can be fulfilled if $`a=0`$ that leads to Eqs (1)-(2). Unlike metric theories of gravitation, in PGTG terms of $`L_G`$ quadratic in the curvature tensor do not lead to high derivatives in GCFE. Parameter $`\alpha `$ in Eqs. (II)–(II)is defined as $`\alpha =\frac{1}{3}\left(16\pi \right)^2fM_p^4`$. In our previous papers the parameter $`\beta =\alpha `$ was used.:
$`{\displaystyle \frac{k}{R^2}}+\left\{{\displaystyle \frac{d}{dt}}\mathrm{ln}\left[R\sqrt{\left|1+\alpha \left(\rho 3p\right)\right|}\right]\right\}^2`$
$`={\displaystyle \frac{8\pi }{3M_p^2}}{\displaystyle \frac{\rho +\frac{\alpha }{4}\left(\rho 3p\right)^2}{1+\alpha \left(\rho 3p\right)}},`$ (1)
$`{\displaystyle \frac{\left[\dot{R}+R\left(\mathrm{ln}\sqrt{\left|1+\alpha \left(\rho 3p\right)\right|}\right)^{}\right]^{}}{R}}`$
$`={\displaystyle \frac{4\pi }{3M_p^2}}{\displaystyle \frac{\rho +3p\frac{\alpha }{2}\left(\rho 3p\right)^2}{1+\alpha \left(\rho 3p\right)}},`$ (2)
where $`k=+1,0,1`$ for closed, flat, open models respectively, $`\alpha `$ is indefinite parameter with inverse dimension of energy density, $`M_p`$ is Planckian mass, a dot denotes differentiation with respect to time. (The system of units with $`\mathrm{}=c=1`$ is used). Note that Eqs.(II)–(II) are valid also in the frame of the most general gauge theory of gravitation – metric-affine theory mc25 ; mc26 . From Eqs. (II)-(II) follows the conservation law in usual form
$$\dot{\rho }+3H\left(\rho +p\right)=0,$$
(3)
where $`H=\frac{\dot{R}}{R}`$ is the Hubble parameter. Besides cosmological equations (II)–(II) gravitational equations of PGTG lead to the following relation for the torsion function $`S`$
$$S=\frac{1}{4}\frac{d}{dt}\mathrm{ln}\left|1+\alpha (\rho 3p)\right|.$$
(4)
Before we start discussing inflationary cosmology on a base of GCFE (II)-(II), let us make several remarks about these equations. At $`\alpha 0`$ the GCFE (II)-(II) are transformed into Friedmann cosmological equations of GR. The value of $`\alpha ^1`$ determines the scale of extremely high energy densities. Solutions of GCFE (II)–(II) coincide practically with corresponding solutions of GR if the energy density is small $`\left|\alpha (\rho 3p)\right|1`$ ($`p\frac{1}{3}\rho `$). The difference between GR and PGTG can be essential at extremely high energy densities $`\left|\alpha (\rho 3p)\right|1`$. Ultrarelativistic matter ($`p=\frac{1}{3}\rho `$) and gravitating vacuum ($`p=\rho `$) with constant energy density are two exceptional systems because Eqs. (II)–(II) are identical to Friedmann cosmological equations of GR in these cases independently on values of energy density and $`S=0`$. The behaviour of solutions of Eqs. (II)–(II) depends essentially on the following restriction on equation of state of gravitating matter at extreme conditions: $`p>\frac{1}{3}\rho `$ or $`p<\frac{1}{3}\rho `$. Because of strong nucleon interaction we have for nuclear matter $`p>\frac{1}{3}\rho `$ zel ; we will put that this restriction is valid also for gravitating matter at extreme conditions, and the scale of extremely high energy densities determined by $`\alpha ^1`$ surpasses nuclear density. From physical point of view, we can suppose that the value of $`\alpha ^1`$ is smaller than the Planckian energy density, but it will be shown later, the behaviour of solutions of GCFE does not depend on this assumption.
The GCFE lead to restrictions on admissible values of energy density. In fact, if energy density $`\rho `$ is positive and $`\alpha >0`$, from Eq.(II) in the case $`k=+1`$, $`0`$ follows the relation:
$$Z1+\alpha \left(\rho 3p\right)0.$$
(5)
The condition (5) is valid not only for closed and flat models, but also for cosmological models of open type ($`k=1`$) (see below). In the case of models filled by gravitating matter with equation of state $`p=p(\rho )`$, it is easy to obtain the solution of the system of Eqs. (II) and (3) in quadratures mc25\_ ; these solutions are regular in metrics, Hubble parameter, its time derivative and have bouncing character. The transition from compression to expansion takes place by reaching limiting energy density defined by the following condition $`Z=0`$. At first the conclusion on possible existence of limiting energy density for gravitating systems, close by which the gravitational interaction has the character of repulsion, was obtained in Ref.mc14 . Later the idea on limiting energy density as ”universal law of the nature” was discussed by M.A. Markov mark , and the value of limiting energy density was postulated to be equal to the Planckian one. According to (5) the value of limiting energy density in our case depends on parameter $`\alpha `$ and equation of state for gravitating matter at extreme conditions and can be essentially smaller than the Planckian one <sup>3</sup><sup>3</sup>3The GCFE lead to limiting energy density also in the case $`\alpha <0`$ if at extreme conditions $`p<\frac{1}{3}\rho `$ mc14 ..
## III III. Inflationary cosmological models and its properties
Now by using GCFE (II)–(II) we will study homogeneous isotropic models filled by interacting scalar field $`\varphi `$ minimally coupled with gravitation and gravitating matter with equation of state in general form $`p_m=p_m(\rho _m)`$. (The generalization for the case with several scalar fields can be made directly). Then the energy density $`\rho `$ and pressure $`p`$ take the form
$$\rho =\frac{1}{2}\dot{\varphi }^2+V+\rho _m(\rho >0),p=\frac{1}{2}\dot{\varphi }^2V+p_m,$$
(6)
where scalar field potential $`V=V(\varphi ,\rho _m)`$ includes the interaction between scalar field and gravitating matter. In the most important particular case of radiation (ultrarelativistic matter) the expressions of $`V(\varphi ,\rho _m)`$ can be obtained by taking into account temperature corrections for given scalar field potentials mc27 and the following relation for energy density $`\rho _mT^4`$. Because the form of equation of state for gravitating matter changes by cosmological evolution and we do not know explicit form of scalar field potential, our analysis will be made without its concretization. Our main aim will be to investigate properties of inflationary cosmological solutions for early Universe. By using scalar field equation in homogeneous isotropic space
$$\ddot{\varphi }+3H\dot{\varphi }=\frac{V}{\varphi }$$
(7)
we obtain from Eqs. (3), (6), (7) the conservation law for gravitating matter
$$\dot{\rho }_m\left(1+\frac{V}{\rho _m}\right)+3H\left(\rho _m+p_m\right)=0.$$
(8)
By using Eqs.(6)–(8) the GCFE (II)–(II) can be transformed to the following form
$$\begin{array}{c}\left\{H\left(Z+3\alpha \dot{\varphi }^2+\frac{3\alpha }{2}Y\right)+3\alpha \frac{V}{\varphi }\dot{\varphi }\right\}^2\hfill \\ \hfill +\frac{k}{R^2}Z^2=\frac{8\pi }{3M_p^2}\left[\rho _m+\frac{1}{2}\dot{\varphi }^2+V+\frac{1}{4}\alpha \left(4V\dot{\varphi }^2+\rho _m3p_m\right)^2\right]Z,\end{array}$$
(9)
$$\begin{array}{c}\dot{H}[Z+3\alpha (\dot{\varphi }^2+\frac{1}{2}Y)]Z+H^2\{[Z15\alpha \dot{\varphi }^2+\frac{3\alpha }{2}Y\hfill \\ \hfill \frac{9\alpha }{2}\frac{Y}{1+\frac{V}{\rho _m}}(1+\frac{dp_m}{d\rho _m}+\frac{\rho _m+p_m}{1+\frac{V}{\rho _m}}\frac{^2V}{\rho _m^2})\frac{9\alpha }{2}\frac{\left(\rho _m+p_m\right)^2}{\left(1+\frac{V}{\rho _m}\right)^2}\frac{d^2}{d\rho _m^2}(3p_m4V)]Z\\ \hfill 18\alpha ^2(\dot{\varphi }^2+\frac{1}{2}Y)^2\}12\alpha H\dot{\varphi }\{3\alpha \frac{V}{\varphi }(\dot{\varphi }^2+\frac{1}{2}Y)\\ \hfill +[\frac{V}{\varphi }+\frac{9}{8}\frac{\rho _m+p_m}{\left(1+\frac{V}{\rho _m}\right)^2}\frac{^2V}{\varphi \rho _m}(1+\frac{1}{3}\frac{d}{d\rho _m}(p_m+2V))]Z\}\\ \hfill +3\alpha \left[\frac{^2V}{\varphi ^2}\dot{\varphi }^2\left(\frac{V}{\varphi }\right)^2\right]Z18\alpha ^2\left(\frac{V}{\varphi }\right)^2\dot{\varphi }^2\\ \hfill =\frac{8\pi }{3M_p^2}\left[V\dot{\varphi }^2\frac{1}{2}\left(\rho _m+3p_m\right)+\frac{1}{4}\alpha \left(4V\dot{\varphi }^2+\rho _m3p_m\right)^2\right]Z,\end{array}$$
(10)
where $`Z=1+\alpha \left(4V\dot{\varphi }^2+\rho _m3p_m\right)`$ and $`Y=\frac{\rho _m+p_m}{1+{\scriptscriptstyle \frac{V}{\rho _m}}}\left(\frac{d}{d\rho _m}\left(3p_m4V\right)1\right)`$. The formula (4) takes the form
$$S=\frac{3\alpha }{2Z}\left[H\left(\dot{\varphi }^2+\frac{1}{2}Y\right)+\frac{V}{\varphi }\dot{\varphi }\right].$$
(11)
The relation (5) determines the following restriction on admissible values of variables for scalar field and gravitating matter:
$$\dot{\varphi }^24V+\alpha ^1+\rho _m3p_m.$$
(12)
Now let us introduce the 3-dimensional space $`P`$ with axes $`(\varphi ,\dot{\varphi },\rho _m)`$. The domain of admissible values of scalar field $`\varphi `$, time derivative $`\dot{\varphi }`$ and energy density $`\rho _m`$ in space $`P`$ determined by (12) is limited by bound $`L`$ defined as
$$Z=0\text{or}\dot{\varphi }=\pm \left(4V+\alpha ^1+\rho _m3p_m\right)^{\frac{1}{2}}.$$
(13)
From Eq. (9) the Hubble parameter on the bound $`L`$ is equal to
$$H_L=\frac{{\displaystyle \frac{V}{\varphi }}\dot{\varphi }}{\dot{\varphi }^2+{\displaystyle \frac{1}{2}}Y}.$$
(14)
Let us consider the most important general properties of cosmological solutions of GCFE (9)–(10). At first note, by given initial conditions for variables ($`\varphi `$, $`\dot{\varphi },\rho _m`$) and also in the case $`k=\pm 1`$ for $`R`$ there are two different solutions corresponding to two values of the Hubble parameter following from Eq. (9):
$$H_\pm =\frac{3\alpha {\displaystyle \frac{V}{\varphi }}\dot{\varphi }\pm \sqrt{D}}{Z+3\alpha [\dot{\varphi }^2+{\displaystyle \frac{1}{2}}Y]},$$
(15)
where
$$\begin{array}{c}D=\frac{8\pi }{3M_p^2}[\rho _m+\frac{1}{2}\dot{\varphi }^2+V\hfill \\ \hfill +\frac{1}{4}\alpha (4V\dot{\varphi }^2+\rho _m3p_m)^2]Z\frac{k}{R^2}Z^2.\end{array}$$
(16)
Obviously, the expression (15) for $`H_\pm `$ will be regular, if $`Y0`$. This relation is valid, in particular, for all models filled by gravitating matter with $`p_m\frac{\rho _m}{3}`$ and scalar fields with potentials applying in chaotic inflation theory. Unlike GR, the values of $`H_+`$ and $`H_{}`$ in GTG are sign-variable and, hence, both solutions corresponding to $`H_+`$ and $`H_{}`$ can describe the expansion as well as the compression in dependence on its sign. Below we will call solutions of GCFE corresponding to $`H_+`$ and $`H_{}`$ as $`H_+`$-solutions and $`H_{}`$-solutions respectively. Note that at asymptotics like GR the sign of $`H_{}`$ is negative and the sign of $`H_+`$ is positive. In points of bound $`L`$ we have $`D=0`$, $`H_+=H_{}`$ and the Hubble parameter is determined by (14). If initial conditions correspond to asymptotics of $`H_{}`$-solution, then unlike GR by compression stage the derivative $`\dot{\varphi }`$ does not diverge and by reaching the bound $`L`$ the transition from $`H_{}`$-solution to $`H_+`$-solution takes place. In fact, by using the following formula for $`H_\pm `$-solutions
$$\dot{Z}=6\alpha \frac{{\displaystyle \frac{V}{\varphi }}\dot{\varphi }Z\pm \sqrt{D}\left(\dot{\varphi }^2+{\displaystyle \frac{1}{2}}Y\right)}{Z+3\alpha \left(\dot{\varphi }^2+{\displaystyle \frac{1}{2}}Y\right)}$$
(17)
it is easy to show that
$$\underset{Z0}{lim}\dot{H}_\pm =\dot{H}_L+\frac{2\pi }{3\alpha M_p^2}\frac{\rho _m+p_m+\dot{\varphi }^2}{\dot{\varphi }^2+\frac{1}{2}Y}.$$
(18)
From (18) follows that in points of bound $`L`$ the derivatives $`\dot{H}_+`$ and $`\dot{H}_{}`$ are equal and its values do not depend on the model type, as a result we have the smooth transition from $`H_{}`$-solution to $`H_+`$-solution on bound $`L`$, and corresponding cosmological solutions for all types models are regular in metrics, Hubble parameter and its time derivative. Note, that in points of bound $`L`$ conditions of uniqueness of solutions of Eqs. (9)-(10) are not fulfilled, as a result there are specific solutions mc3 ; mc15 , trajectories of which are situated on the bound $`L`$ and have with $`H_\pm `$-solutions common points, $`H_{}`$–solutions reach the bound $`L`$ and $`H_+`$-solutions originate from them, the surface $`Z=0`$ is envelope in space $`P`$ for cosmological solutions <sup>4</sup><sup>4</sup>4From mathematical point of view the appearance of specific solutions is connected with the fact, that the Eqs. (1)–(2) were multiplied by $`Z^2`$ when its transformation to the (9)–(10) was performed. Specific solutions can be excluded by corresponding transformation of Eq. (10).. According to Eq. (17) and (4) the function $`S`$ has the following asymptotics for $`H_+`$\- and $`H_{}`$-solutions at $`Z0`$:
$$\underset{Z0}{lim}S=\frac{1}{4}\underset{Z0}{lim}\frac{\dot{Z}}{Z}Z^{\frac{1}{2}}.$$
(19)
Unlike flat and open models, for which $`H_+=H_{}`$ only in points of bound $`L`$ and regular inflationary models include $`H_+`$– and $`H_{}`$–solutions reaching bound $`L`$, in the case of closed models the regular transition from $`H_{}`$-solution to $`H_+`$-solution is possible without reaching the bound $`L`$. It is because by certain value of $`R`$ according to (15)–(16) we have $`H_+=H_{}`$ in the case $`Z0`$. Such models are regular also in torsion . Regular inflationary solution of such type was considered in Ref.mc6 .
All discussed cosmological solutions have bouncing character, but in the presence of scalar fields a bounce takes place not in points of bound $`L`$ ($`Z=0`$). In order to study the behaviour of cosmological models at a bounce, let us analyze extreme points for the scale factor $`R(t)`$: $`R_0=R(0)`$, $`H_0=H(0)=0`$. (This means that in the case of $`H_+`$–solutions $`H_{+0}=0`$ and in the case of $`H_{}`$–solutions $`H_0=0`$). Denoting values of quantities at $`t=0`$ by means of index ”0”, we obtain from (9)–(10):
$$\frac{k}{R_0^2}Z_0^2+9\alpha ^2(\frac{V}{\varphi })_0^2\dot{\varphi }_0^2=\frac{8\pi }{3M_p^2}\left[\rho _{m0}+\frac{1}{2}\dot{\varphi }_0^2+V_0+\frac{1}{4}\alpha \left(4V_0\dot{\varphi }_0^2+\rho _{m0}3p_{m0}\right)^2\right]Z_0,$$
(20)
$$\begin{array}{c}\dot{H}_0=\{\frac{8\pi }{3M_p^2}[V_0\dot{\varphi }_0^2\frac{1}{2}(\rho _{m0}+3p_{m0})+\frac{1}{4}\alpha (4V_0\dot{\varphi }_0^2+\rho _{m0}3p_{m0})^2]\hfill \\ \hfill 3\alpha [(\frac{^2V}{\varphi ^2})_0\dot{\varphi }_0^2(\frac{V}{\varphi })_0^2]+18\alpha ^2(\frac{V}{\varphi })_0^2\dot{\varphi }_0^2Z_0^1\left\}\right\{Z_0+3\alpha [\dot{\varphi }_0^2+\frac{1}{2}Y_0]\}^1,\end{array}$$
(21)
where $`Z_0=1+\alpha \left(4V_0\dot{\varphi }_0^2+\rho _{m0}3p_{m0}\right)`$. A bounce point is described by Eq. (20), if the value of $`\dot{H}_0`$ is positive. By using Eq.(20) we can rewrite the expression of $`\dot{H}_0`$ in the form
$`\dot{H}_0=\{{\displaystyle \frac{8\pi }{M_p^2}}[V_0+{\displaystyle \frac{1}{2}}(\rho _{m0}p_{m0})+{\displaystyle \frac{1}{4}}\alpha (4V_0\dot{\varphi }_0^2+\rho _{m0}3p_{m0})^2]`$
$`3\alpha [({\displaystyle \frac{^2V}{\varphi ^2}})_0\dot{\varphi }_0^2({\displaystyle \frac{V}{\varphi }})_0^2]{\displaystyle \frac{2k}{R_0^2}}Z_0\left\}\right\{Z_0+3\alpha [\dot{\varphi }_0^2+{\displaystyle \frac{1}{2}}Y_0]\}^1.`$ (22)
We see from (22) unlike GR the presence of gravitating matter satisfying the energy dominance condition ($`p_m\rho _m`$) does not prevent from the bounce realization<sup>5</sup><sup>5</sup>5In GR a bounce is possible only in closed models if the following condition $`V_0\dot{\varphi }_0^2\frac{1}{2}(\rho _{m0}+3p_{mo})>0`$ takes place.. Eq. (20) determines in space $`P`$ extremum surfaces depending on the value of $`\alpha `$ and in the case of closed and open models also parametrically on the scale factor $`R_0`$. In the case of various scalar field potentials applying in inflationary cosmology the value of $`\dot{H}_{+0}`$ or $`\dot{H}_0`$ is positive on the greatest part of extremum surfaces, which can be called ”bounce surfaces”<sup>6</sup><sup>6</sup>6If $`\alpha ^1M_p^4`$ the derivative $`\dot{H}_0`$ is negative in the neighbourhood of origin of coordinates in space $`P`$ that leads to appearance of oscillating solutions of GCFEmc28 .. By giving concrete form of potential $`V`$ and choosing values of $`R_0`$, $`\varphi _0`$, $`\dot{\varphi }_0`$ and $`\rho _{m0}`$ at a bounce, we can obtain numerically particular bouncing solutions of GCFE for various values of parameter $`\alpha `$. Like GR, if initial value of scalar field at the beginning of cosmological expansion is not small ($`\varphi 1M_p`$) mc27 , corresponding solution is inflationary cosmological solution containing in addition to inflationary stage also compression stage, transition stage from compression to expansion and a stage after inflation with oscillating regime for scalar field. For given scalar field potential properties of regular inflationary cosmological solutions depend on initial conditions at a bounce and parameter $`\alpha `$. Numerical analysis of such solutions in the case of simplest scalar field potentials applying in chaotic inflation was carried out in Ref. mc29 . Note that by equal initial conditions characteristics of inflationary stage in developing thery coincide with that of GR mc29 .
The analysis of GCFE shows, that some properties of cosmological solutions depend essentially on parameter $`\alpha `$, i.e. on the scale of extremely high energy densities. From physical point of view interesting results can be obtained, if the value of $`\alpha ^1`$ is much less than the Planckian energy densitymc6 , i.e. in the case of large values of parameter $`\alpha `$ (by imposing $`M_p=1`$). In order to investigate cosmological solutions at the beginning of cosmological expansion in this case, let us consider the GCFE by supposing that
$`\left|\alpha \left(4V\dot{\varphi }^2+\rho _m3p_m\right)\right|1,`$
$`\rho _m+{\displaystyle \frac{1}{2}}\dot{\varphi }^2+V\alpha \left(4V\dot{\varphi }^2+\rho _m3p_m\right)^2.`$ (23)
Note that the second condition (III) does not exclude that ultrarelativistic matter energy density can dominate at a bounce. We obtain:
$$\begin{array}{c}\frac{k}{R^2}(4V\dot{\varphi }^2+\rho _m3p_m)^2+\{H[4V+2\dot{\varphi }^2+\rho _m3p_m+\frac{3}{2}Y]\hfill \\ \hfill +3\frac{V}{\varphi }\dot{\varphi }\}^2=\frac{2\pi }{3M_p^2}(4V\dot{\varphi }^2+\rho _m3p_m)^3,\end{array}$$
(24)
$$\begin{array}{c}\dot{H}(4V+2\dot{\varphi }^2+\rho _m3p_m+\frac{3}{2}Y)(4V\dot{\varphi }^2+\rho _m3p_m)+H^2\{[4V16\dot{\varphi }^2+\rho _m3p_m\hfill \\ \hfill +\frac{3}{2}Y\frac{9}{2}\frac{Y}{1+\frac{V}{\rho _m}}(1+\frac{dp_m}{d\rho _m}+\frac{\rho _m+p_m}{1+\frac{V}{\rho _m}}\frac{^2V}{\rho _m^2})\frac{9}{2}\frac{\left(\rho _m+p_m\right)^2}{\left(1+\frac{V}{\rho _m}\right)^2}\frac{d^2}{d\rho _m^2}(3p_m4V)]\\ \hfill \times (4V\dot{\varphi }^2+\rho _m3p_m)18(\dot{\varphi }^2+\frac{1}{2}Y)^2\}\\ \hfill 12H\dot{\varphi }\{[\frac{V}{\varphi }+\frac{9}{8}\frac{\rho _m+p_m}{\left(1+\frac{V}{\rho _m}\right)^2}\frac{^2V}{\varphi \rho _m}(1+\frac{1}{3}\frac{d}{d\rho _m}(p_m+2V))]\\ \hfill \times (4V\dot{\varphi }^2+\rho _m3p_m)+3\frac{V}{\varphi }(\dot{\varphi }^2+\frac{1}{2}Y)\}\\ \hfill +3\left[\frac{^2V}{\varphi ^2}\dot{\varphi }^2\left(\frac{V}{\varphi }\right)^2\right]\left(4V\dot{\varphi }^2+\rho _m3p_m\right)\\ \hfill 18\left(\frac{V}{\varphi }\right)^2\dot{\varphi }^2=\frac{2\pi }{3M_p^2}\left(4V\dot{\varphi }^2+\rho _m3p_m\right)^3.\end{array}$$
(25)
According to Eq. (24) the Hubble parameter in considered approximation is equal to
$$\begin{array}{c}H_\pm =\left[4V+2\dot{\varphi }^2+\rho _m3p_m+\frac{3}{2}Y\right]^1\hfill \\ \hfill \times \left[3\frac{V}{\varphi }\dot{\varphi }\pm \left|4V\dot{\varphi }^2+\rho _m3p_m\right|\sqrt{\frac{2\pi }{3M_p^2}\left(4V\dot{\varphi }^2+\rho _m3p_m\right)\frac{k}{R^2}}\right]\end{array}$$
(26)
and extreme points of the scale factor are determined by the following condition
$`{\displaystyle \frac{k}{R_0^2}}+9\left[{\displaystyle \frac{\left({\displaystyle \frac{V}{\varphi }}\right)_0\dot{\varphi }_0}{4V_0\dot{\varphi }_0^2+\rho _{m0}3p_{m0}}}\right]^2`$
$`={\displaystyle \frac{2\pi }{3M_p^2}}\left(4V_0\dot{\varphi }_0^2+\rho _{m0}3p_{m0}\right).`$ (27)
From Eq. (25) the time derivative of the Hubble parameter at extreme points is
$$\begin{array}{c}\dot{H}_0=\{\frac{2\pi }{3M_p^2}(4V_0\dot{\varphi }_0^2+\rho _{m0}3p_{m0})^2\hfill \\ \hfill 3[(\frac{^2V}{\varphi ^2})_0\dot{\varphi }_0^2(\frac{V}{\varphi })_0^2]+\frac{18(\frac{V}{\varphi })_0^2\dot{\varphi }^2}{4V_0\dot{\varphi }_0^2+\rho _{m0}3p_{m0}}\}\\ \hfill \times \left[4V_0+2\dot{\varphi }_0^2+\rho _{m0}3p_{m0}+\frac{3}{2}Y_0\right]^1.\end{array}$$
(28)
From Eqs.(24)–(28) follows that in approximation (III) the dynamics of inflationary cosmological models at a bounce does not depend on parameter $`\alpha `$. In the case of models containing at a bounce ultrarelativistic matter and scalar fields, the dynamics in considered approximation does not depend on ultrarelativistic matter, if interaction terms in scalar field potential can be neglected.
The analysis given in this Section shows that existence of limiting bound $`L`$ and bounce surface in space $`P`$ ensures regular character of inflationary cosmological solutions in metrics, Hubble parameter and its time derivative. Characteristic feature of inflationary cosmological models is the presence of the stage of regular transition from compression to expansion. The duration of this stage is several order smaller than duration of inflationary stage. If we take into account that duration of inflationary stage in chaotic inflation is extremely small mc27 , we can tell that our inflationary models correspond to regular Big Bang or Big Bounce.
## IV IV. Particular regular inflationary cosmological model
As illustration of discussed theory regular cosmological models in the simplest particular case will be considered in this Section. We will consider models including noninteracting scalar field with potential $`V(\varphi )`$ and ultrarelativistic matter ($`p_m=\frac{1}{3}\rho _m`$)mc2 . In this case the bound $`L`$ in space $`P`$ is reduced to two cylindric surfaces $`\dot{\varphi }=\pm \left(4V+\alpha ^1\right)^{\frac{1}{2}}`$. Bounce surface is reduced also to cylindric surfaces in the case under consideration, when the scale of extremely high energy densities is much smaller than the Planckian energy density (see (III)). In connection with this we will consider instead of space $`P`$ the plane of variables ($`\varphi `$, $`\dot{\varphi }`$) and intersections of the bound $`L`$ and bounce surface with this plane. We have in this plane two bound $`L_\pm `$–curves and in the case of flat models two bounce curves $`B_1`$ and $`B_2`$ determined by equation <sup>7</sup><sup>7</sup>7The neighbourhood of origin of coordinates is not considered in this approximation, the behavior of extremum curves near origin of coordinates was examined in Ref. mc28 ; mc29 .
$$4V_0\dot{\varphi }_0^2=3\left(\frac{M_p^2}{2\pi }V_{0}^{}{}_{}{}^{2}\dot{\varphi }_0^2\right)^{\frac{1}{3}}.$$
Each of two curves $`B_{1,2}`$ contains two parts corresponding to vanishing of $`H_+`$ or $`H_{}`$ and denoting by ($`B_{1+}`$, $`B_{2+}`$) and ($`B_1`$, $`B_2`$) respectively. If $`V^{}`$ is positive (negative) in quadrants 1 and 4 (2 and 3) on the plane ($`\varphi `$, $`\dot{\varphi }`$), the bounce will take place in points of bounce curves $`B_{1+}`$ and $`B_{2+}`$ ($`B_1`$ and $`B_2`$) in quadrants 1 and 3 (2 and 4) for $`H_+`$-solutions ($`H_{}`$-solutions) (see Fig. 1).
To analyze flat bouncing models we have to take into account that besides regions lying between curves $`L_\pm `$ and corresponding bounce curves the sign of values $`H_+`$ and $`H_{}`$ for applying potentials is normal: $`H_+>0`$, $`H_{}<0`$. The Hubble parameter $`H_+`$ is negative in regions between curves ($`L_+`$ and $`B_{1+}`$), ($`L_{}`$ and $`B_{2+}`$), and the value of $`H_{}`$ is positive in regions between curves ($`L_+`$ and $`B_1`$), ($`L_{}`$ and $`B_2`$). As it was noted above any cosmological solution has to contain both $`H_{}`$\- and $`H_+`$-solution and regular transition from $`H_{}`$-solution to $`H_+`$-solution takes place in points of $`L_\pm `$ where $`H_+=H_{}`$. In the case of open and closed models Eq.(III) determines 1-parametric family of bounce curves with parameter $`R_0`$. Bounce curves of closed models are situated in region between two bounce curves $`B_1`$ and $`B_2`$ of flat models, and in the case of open models bounce curves are situated in two regions between the curves: $`L_+`$ and $`B_1`$, $`L_{}`$ and $`B_2`$. In general case, when approximation (III) is not valid, bounce surface in space $`P`$ of cosmological models including scalar field and ultrarelativistic matter determined by Eq.(20) in space $`P`$ depends on parameter $`\alpha `$ and it is not more cylindric surface. The situation concerning cosmological solutions of Eqs. (9)–(10) does not change.
Below particular bouncing cosmological inflationary solution for flat model by using scalar field potential in the form $`V=\frac{1}{2}m^2\varphi ^2`$ ($`m=10^6M_p`$)is given. The solution was obtained by numerical integration of Eqs. (7), (10) and by choosing in accordance with Eq.(20) the following initial conditions at a bounce: $`\varphi _0=\sqrt{2}\mathrm{\hspace{0.17em}10}^3M_p`$, $`\dot{\varphi }_0=\sqrt{3.96757V_0}`$ ($`\alpha =10^{14}M_p^4`$); radiation energy density is negligibly small, initial value of $`R_0`$ can be arbitrary. The dynamics of the Hubble parameter and scalar field is presented for different stages of obtained bouncing solution in Figs. 24 (by choosing $`M_p=1`$).
The transition stage from compression to expansion (Fig. 2) is essentially asymmetric with respect to the point $`t=0`$ because of $`\dot{\varphi }_00`$. In course of transition stage the Hubble parameter changes from maximum in module negative value at the end of compression stage to maximum positive value at the beginning of expansion stage. The scalar field changes linearly, the derivative $`\dot{\varphi }`$ grows at first being positive to maximum value $`\dot{\varphi }\dot{\varphi }_0`$ and then the value of $`\dot{\varphi }`$ decreases and becomes negative. In the case of presence of radiation, which can give the main contribution to energy density at a bounce, the duration of transition stage changes. Quasi-de-Sitter inflationary stage and quasi-de-Sitter compression stage are presented in Fig. 3. As was noted above, characteristics of inflationary stage do not depend practically on parameter $`\alpha `$ and coincide with that of GR. The amplitude and frequency of oscillating scalar field after inflation (Fig. 4) are close to that of GR, however, the behaviour of the Hubble parameter after inflation in considering case with large value of parameter $`\alpha `$ is essentially noneinsteinian, at first the Hubble parameter oscillates near the value $`H=0`$, and later the Hubble parameter becomes positive and decreases with the time like in GR. Before quasi-de-Sitter compression stage there are also oscillations of the Hubble parameter and scalar field not presented in Figs. 24. Ultrarelativistic matter, which could dominate at a bounce has negligibly small energy densities at quasi-de Sitter stages. At the same time the gravitating matter could be at compression stage in more realistic bouncing models. As it follows from our consideration regular character of such inflationary cosmological models has to be ensured by cosmological equations of PGTG.
The interaction between scalar fields and radiation leads to quantitative corrections of considered cosmological models. In accordance with Eq. (8) temperature corrections for scalar field potentials change the connection between scale factor $`R`$ and radiation energy density and can be also essential for more late stages of cosmological evolution, when energy densities are sufficiently small and consequences of PGTG and GR coincide. In particular, these corrections can give additional contribution to the effect of acceleration of cosmological expansion (if relic scalar fields exist).
Note that bouncing character have solutions not only in classical region, where scalar field potential, kinetic energy density of scalar field and energy density of gravitating matter do not exceed the Planckian energy density, but also in regions, where classical restrictions are not fulfilled and according to accepted opinion quantum gravitational effects can be essential.
## V V. Conclusion
As we see, the applying of gauge approach to gravitational interaction permits to build consequent field theoretical scheme in the frame of 4-dimensional physical space-time, which is free of principal difficulties of GR by description of early Universe. Satisfying the correspondence principle with GR in the case of gravitating systems with rather small energy densities, generalized cosmological Friedmann equations lead to conclusion, that at extreme conditions gravitational interaction has the repulsion character in the case of usual gravitating systems with positive values of energy density satisfying the energy dominance condition. This means, there is not necessity to refuse fundamental physical requirement of energy density positivity for physical matter. The solution of PCS is obtained by classical description of gravitational field without quantum gravitational corrections. Moreover, from the point of view of developing approach, the Planckian era could be absent by evolution of our Universe. Unlike loop quantum cosmology, the Universe is macro-object at all stages of its evolution including the transition stage from compression to expansion. To build a realistic cosmological model we have to know the content and properties of gravitating matter at different stages of its evolution. It is of principal interest to investigate physical processes at the beginning of cosmological expansion depending on limiting energy density and limiting temperature, and to obtain observable physical consequences depending on parameter $`\alpha `$. From physical point of view, it is interesting also the building of cosmological models by breakdown of its homogeneity and isotropy. |
warning/0506/hep-ph0506102.html | ar5iv | text | # OUTP–0505P TUM-HEP–591/05 hep–ph/0506102 Neutrinoless Double Beta Decay and Future Neutrino Oscillation Precision Experiments
## 1 Introduction
The last few years have seen tremendous progress in our understanding of the neutrino sector. The solar neutrino deficit is now known to be certainly due to neutrino flavor oscillations with the best–fit oscillation parameters given as $`\mathrm{\Delta }m_{}^2\mathrm{\Delta }m_{21}^2=8.0\times 10^5`$ eV<sup>2</sup> and $`\mathrm{sin}^2\theta _{}\mathrm{sin}^2\theta _{12}=0.31`$ . The atmospheric neutrino deficit can also be ascribed to flavor oscillations with a good degree of confidence with the best–fit oscillation parameters as $`\mathrm{\Delta }m_{\mathrm{atm}}^2|\mathrm{\Delta }m_{31}^2|=2.1\times 10^3`$ eV<sup>2</sup> and $`\mathrm{sin}^22\theta _{\mathrm{atm}}\mathrm{sin}^22\theta _{23}=1`$ . The upper limit on the third mixing angle $`\theta _{13}`$ is mainly determined by the reactor neutrino data , which when combined with the information obtained from the solar and atmospheric neutrino experiments gives a bound of $`\mathrm{sin}^2\theta _{13}<0.044`$ at $`3\sigma `$ .
Despite all these impressive achievements, there still remains a lot to be learned. The arguably most fundamental question, namely whether neutrinos are Dirac or Majorana particles, remains still to be answered. If neutrinos are Majorana particles, we have nine physical parameters describing the $`3\times 3`$ light neutrino mass matrix and determining them is one of the ultimate goals of neutrino physics. With the two $`\mathrm{\Delta }m^2`$ and two mixing angles – $`\theta _{12}`$ and $`\theta _{23}`$ — known, there still remain five parameters to be tracked down. To be more precise, we still lack knowledge about the following points:
* Is there $`CP`$ violation in the lepton sector in analogy with that in the quark sector?
* Is the atmospheric neutrino mixing angle $`\theta _{23}`$ exactly maximal? If not, does it lie above or below $`\pi /4`$?
* Is the third mixing angle $`\theta _{13}`$ exactly zero?
* How small is the absolute neutrino mass scale?
* What is the ordering of the neutrino masses, i.e., what is the neutrino mass hierarchy?
Answers to all or some of these question will certainly lead to better understanding of the underlying theory which gives rise to neutrino masses<sup>1</sup><sup>1</sup>1 A rarely mentioned point is that for Majorana neutrinos one can, at least in principle, fully determine the complete mass matrix because it is symmetric. If the neutrinos are Dirac particles like quarks and charged leptons, their mass matrix is in general not symmetric and therefore parametrized by fifteen free parameters. However, since we only see left–handed weak currents, the “right–handed” parts of the neutrino parameters are out of reach and hence not observable. .
Neutrinoless double beta decay (0$`\nu \beta \beta `$) is expected to be of crucial importance in answering some of the questions raised above. At the same time, answers to some the questions obtained from elsewhere will help in interpreting a positive or negative signal in the next generation 0$`\nu \beta \beta `$ experiments. Goal of these experiments is the observation of the process
$$(A,Z)(A,Z2)+2e^{}.$$
The effective mass which will be extracted or bounded in a 0$`\nu \beta \beta `$ experiment is given by the following coherent sum:
$`m=\left|{\displaystyle \underset{i}{}}m_iU_{ei}^2\right|,`$ (1)
where $`m_i`$ is the mass of the $`i^{th}`$ neutrino mass state, the sum is over all the light neutrino mass states and $`U_{ei}`$ are the matrix elements of the Pontecorvo–Maki–Nakagawa–Sakata (PMNS) neutrino mixing matrix . We parametrize it as
$`U=\left(\begin{array}{ccc}c_{12}c_{13}& s_{12}c_{13}& s_{13}\\ s_{12}c_{23}c_{12}s_{23}s_{13}e^{i\delta }& c_{12}c_{23}s_{12}s_{23}s_{13}e^{i\delta }& s_{23}c_{13}e^{i\delta }\\ s_{12}s_{23}c_{12}c_{23}s_{13}e^{i\delta }& c_{12}s_{23}s_{12}c_{23}s_{13}e^{i\delta }& c_{23}c_{13}e^{i\delta }\end{array}\right)\mathrm{diag}(1,e^{i\alpha },e^{i\beta }),`$ (5)
where we have used the usual notations $`c_{ij}=\mathrm{cos}\theta _{ij}`$, $`s_{ij}=\mathrm{sin}\theta _{ij}`$, $`\delta `$ is the Dirac $`CP`$–violation phase, $`\alpha `$ and $`\beta `$ are the two Majorana $`CP`$–violation phases . Thus we can define $`m`$ in terms of the oscillation parameters, the Majorana phases and the neutrino mass scale. This means that $`m`$ depends on 7 out of 9 parameters contained in the neutrino mass matrix. In particular, the effective mass is a function of all unknowns of neutrino physics except for the Dirac phase and $`\theta _{23}`$. This means that a measurement of – or a stringent limit on – 0$`\nu \beta \beta `$ will give us some information on the unknown parameters. Combined with complementary information from other independent types of experiments, this would reinforce our understanding of the neutrino sector.
There is a large list of available analyzes focusing on the connection of the known and unknown neutrino parameters with the effective mass , for a recent review of the theoretical situation see . In this work we want to focus on the mass ordering of the neutrinos, on the neutrino mass scale and on the implications of a very small or vanishing effective mass. The effective mass to be extracted from neutrinoless double beta decay depends crucially on the neutrino mass spectrum. Of special interest are the following three extreme cases:
normal hierarchy (NH): $`|m_3|\sqrt{\mathrm{\Delta }m_\mathrm{A}^2}|m_2|\sqrt{\mathrm{\Delta }m_{}^2}|m_1|,`$ (6)
inverted hierarchy (IH): $`|m_2||m_1|\sqrt{\mathrm{\Delta }m_\mathrm{A}^2}|m_3|,`$ (7)
quasi–degeneracy (QD): $`m_0|m_3||m_2||m_1|\sqrt{\mathrm{\Delta }m_\mathrm{A}^2}.`$ (8)
The order of magnitude of the effective mass in those spectra is $`\sqrt{\mathrm{\Delta }m_{}^2},\sqrt{\mathrm{\Delta }m_\mathrm{A}^2}`$ and $`m_0`$, respectively.
The best current limit on the effective mass is given by measurements of <sup>76</sup>Ge established by the Heidelberg–Moscow collaboration
$`m0.35z\mathrm{eV},`$ (9)
where $`z=𝒪(1)`$ indicates that there is an uncertainty stemming from the nuclear physics involved in calculating the decay width of 0$`\nu \beta \beta `$. Similar results were obtained by the IGEX collaboration . Several new experiments are currently running, under construction or in the planing phase. The NEMO3 and CUORICINO experiments are already taking data and reach sensitivities near the current limits. The next generation of experiments, with projects such as CUORE , MAJORANA , GERDA , EXO , MOON , COBRA , XMASS, DCBA , CANDLES , CAMEO , (for a review see ) will probe values of $`m`$ one order of magnitude below the limit from Eq. (9). Thus we expect $`m`$ to be probed down to $`\sqrt{\mathrm{\Delta }m_\mathrm{A}^2}0.04`$ eV and it would be pertinent to ask if such a measurement could help us learn about the neutrino hierarchy<sup>2</sup><sup>2</sup>2Of course those experiments aim also to put the controversial evidence of part of the Heidelberg–Moscow collaboration to the test..
Maximal mixing in the solar neutrino sector is now disfavored at almost $`6\sigma `$ and this non–maximality of $`\theta _{12}`$ allows for the possibility to distinguish the NH from the IH case. This however depends on the value of the smallest neutrino mass state. We investigate this issue and give the value of the smallest neutrino mass state which allows to distinguish NH from IH. Another aspect, concerning the QD spectrum, is the size of the absolute neutrino mass scale. We give and analyze a simple formula for a limit on $`m_0`$ in terms of $`\theta _{12},\theta _{13}`$ and the limit (or value) of $`m`$. All these issues depend crucially on the uncertainties involved. In addition to the experimental errors, we would have to contend with the theoretical uncertainties coming from unknown nuclear matrix elements and uncertainties in the values of the oscillation parameters. As far the oscillation parameters are concerned, the uncertainties are expected to be reduced sharply by the next generation experiments. This will greatly reduce the expected range of $`m`$ for a given mass spectrum. We give a detailed review of the improvements expected in each of the oscillation parameters and identify the parameters which would still lead to maximum uncertainty in the predicted value of $`m`$ for a given mass spectrum. We take the uncertainties coming from our lack of knowledge of the nuclear matrix elements and explore the chances of determining the neutrino mass hierarchy from a positive 0$`\nu \beta \beta `$ signal in the future. Furthermore, we analyze the consequences of a very small or even vanishing $`m`$ if neutrinos are indeed Majorana particles. Such a small $`m`$ will influence the allowed range of values of the mixing parameters $`\theta _{12}`$, $`\theta _{13}`$, the Majorana phases and the smallest neutrino mass. We compare these implied range of parameter values with their current limits. Finally, we perform a thorough scan of the phenomenologically viable neutrino mass models and give a list of possible neutrino mass matrices out of which the true one might be identified when future precision measurements of neutrino oscillation observables are combined with a measurement of or limit on $`m`$.
We begin by giving in Section 2 a detailed report of the existing range of values for the oscillation parameters. We list down the expected improvements in each of the parameters. In Section 3 we probe the chances of distinguishing the IH from the NH from a future positive signal for neutrinoless double beta decay . We take into account the current and future expected uncertainties on the oscillation parameters and include the role of the uncertainty in the nuclear matrix elements. In Section 4 we consider QD as the true mass spectrum of the neutrinos and study what limits one could set on the absolute/common neutrino mass scale $`m_0`$ from a measured value of $`m`$ and the current and future uncertainties on the oscillation parameters. Section 5 probes the situation one would have if we do not see any signal in the next set of 0$`\nu \beta \beta `$ experiments and yet believe that neutrinos are Majorana particles. In Section 6 we then list viable neutrino mass matrices and discuss how the number can drastically be reduced when measurements of $`m`$ and the other oscillation parameters have been performed. We finally conclude in Section 7.
## 2 Past, Present and Future data
In this Section we shall present the status of the current global neutrino data and the prospects for its future improvement. First we give the current situation.
At the $`3\sigma `$ level, our current knowledge of the solar parameters within $`3\sigma `$ is limited to
$`7.0\times 10^5\text{ eV}^2<\mathrm{\Delta }m_{}^2<9.3\times 10^5\text{ eV}^2,`$ (10)
$`0.24<\mathrm{sin}^2\theta _{12}<0.41.`$ (11)
The atmospheric mass squared difference and mixing angle at $`3\sigma `$ are known within
$`1.3\times 10^3\text{ eV}^2<\mathrm{\Delta }m_\mathrm{A}^2<4.2\times 10^3\text{ eV}^2`$ , (12)
$`0.33<\mathrm{sin}^2\theta _{23}<0.66`$ , (13)
while the mixing angle $`\theta _{13}`$ at $`3\sigma `$ is restricted to lie below the value
$`\mathrm{sin}^2\theta _{13}<0.044.`$ (14)
A very useful parameter related to the mass hierarchy of the neutrinos is the ratio of the solar and atmospheric mass squared differences,
$`R{\displaystyle \frac{\mathrm{\Delta }m_{}^2}{\mathrm{\Delta }m_\mathrm{A}^2}},\text{ with }0.017<R<0.072\mathrm{at}3\sigma .`$ (15)
Regarding the absolute value of the neutrino masses, several approaches exist. The most model–independent one is certainly the direct search for kinematical effects in the energy spectra of beta–decays. The Mainz and Troitsk experiments gave upper limits on the electron neutrino mass of 2.3 eV at 95 $`\%`$ C.L. Cosmological observations imply typically more stringent limits, which however considerably depend on the data set and the priors used in the analysis. Combining the cosmic microwave radiation measurements by the WMAP satellite with data on the large scale structure of the Universe and other data sets, gives limits on the sum of neutrino masses between 0.4 and 2 eV, see for a recent update of the situation.
Finally, there is currently no information on any of the three possible $`CP`$ phases, which take on values between zero and $`2\pi `$.
Let us now discuss the current uncertainties in the values of the parameters and their expected improvement. We can see from Eqs. (10) and (12) that we still have $`53`$% ($`14`$%) uncertainty<sup>3</sup><sup>3</sup>3We define uncertainty of a quantity as the difference of the maximally and minimally allowed value divided by its sum and then multiplied by 100. on the value of $`\mathrm{\Delta }m_{31}^2`$ ($`\mathrm{\Delta }m_{21}^2`$) at $`3\sigma `$. These ranges of allowed values are expected to reduce remarkably in the future: the uncertainty in $`\mathrm{\Delta }m_{21}^2`$ is expected to reduce to $`7`$% at the $`3\sigma `$ level from future measurement of reactor antineutrinos in KamLAND and the value of $`\mathrm{\Delta }m_{31}^2`$ is expected to be known with $`5`$% at the $`3\sigma `$ level from the next generation accelerator experiments involving superbeams in the next ten years .
The uncertainty in the knowledge of $`\mathrm{sin}^2\theta _{23}`$ is at present $`33`$% at $`3\sigma `$. This is expected to reduce to $`20`$% in the next ten years after the results from the next generation superbeam experiments become available . It is clear that the mixing angle $`\theta _{23}`$ plays no role for the study of neutrinoless double beta decay. Nevertheless, we shall see later on in Section 6 that in order to identify the neutrino mass matrix it can be rather important to know “how maximal” $`\theta _{23}`$ actually is. Of particular importance in this respect is the question whether the atmospheric neutrino mixing is exactly maximal because this would point to the presence of an underlying symmetry. This deviation will be known to order $`10\%`$ after the next generation long–baseline experiments . Comparative constraints are expected from future SK atmospheric neutrino data with statistics 50 times the current SK statistics . Another way to state this is that we can identify if $`\theta _{23}\pi /4`$ is of order $`\sqrt{R}`$ or smaller with next generation long–baseline experiments. A deviation from $`\pi /4`$ even smaller (i.e., order $`R`$) can be achieved by dedicated next generation experiments (see also ).
The effective mass in neutrinoless double beta is crucially dependent on the values of the mixing parameters $`\theta _{12}`$ and $`\theta _{13}`$. The current uncertainty on the values of $`\mathrm{sin}^2\theta _{12}`$ is $`29`$% at the $`3\sigma `$ level while $`\mathrm{sin}^2\theta _{13}`$ is completely unknown. The uncertainty on the value of $`\mathrm{sin}^2\theta _{12}`$ and the limit on $`\mathrm{sin}^2\theta _{13}`$ is however expected to reduce in the future. While the conventional beam experiments, MINOS, ICARUS and OPERA are expected to provide moderate improvement on $`\mathrm{sin}^2\theta _{13}`$ , the Double–Chooz reactor experiment in France is expected to reduce the upper limit to $`\mathrm{sin}^2\theta _{13}<0.008`$ at the $`3\sigma `$ level . This upper limit could be improved further to $`\mathrm{sin}^2\theta _{13}<0.0025`$ by the combination of the next generation beam experiments T2K and NO$`\nu `$A , as well as by the second generation reactor experiments . It should be borne in mind that any one of these experiments could even measure a non–zero $`\theta _{13}`$, if the true value of $`\theta _{13}`$ happens to fall within their range of sensitivity.
On the other hand, the improvement expected for $`\mathrm{sin}^2\theta _{12}`$ from the currently running experiments is small . The results from the third and the final Helium phase of the ongoing SNO experiment are expected to reduce the uncertainty on $`\mathrm{sin}^2\theta _{12}`$ to not more than $`21`$% and the KamLAND experiment is not expected to make any significant improvement on it . Doping the Super–Kamiokande (SK) detector with 0.1% of Gadolinium could improve our knowledge on the true value of $`\mathrm{sin}^2\theta _{12}`$ to $`18`$% uncertainty . The proposed/planned future experiments aiming to measure the very low energy $`pp`$ flux coming from the sun are also expected to give a better measurement of the solar mixing angle . However, even with very small experimental uncertainty of only 1%, they are not expected to reduce the uncertainty on $`\mathrm{sin}^2\theta _{12}`$ much better than $`14`$% at $`3\sigma `$ . The only type of experiment that would provide extremely good measurement of the solar mixing angle is a long baseline reactor experiment with its baseline tuned to a Survival Probability minimum , which is given by the condition $`L_{min}1.24E/\mathrm{MeV}\mathrm{eV}^2/\mathrm{\Delta }\mathrm{m}_{21}^2`$. This type of experiment could be used to measure the mixing angle $`\mathrm{sin}^2\theta _{12}`$ down to $`6`$% at $`3\sigma `$ .
The Dirac $`CP`$ phase $`\delta `$ has a faint chance of being measured in the next generation superbeam experiments , provided the true value of $`\theta _{13}`$ is not too small. However, performing an unambiguous measurement, and especially establishing a signal for $`CP`$ violation, would probably need a beta beam facility or a neutrino factory . The two Majorana phases are measurable only in processes which violate lepton number. At present the only such process which seems to be viable experimentally is neutrinoless double beta decay . For detailed analyzes of how to extract information on Majorana phases from 0$`\nu \beta \beta `$ we refer to . The bottom–line of this subject is that for given oscillation parameters the Majorana phases should not be too close to $`0`$ or $`\pi /2`$ and the uncertainty on both the experimental and theoretical side (i.e., the nuclear matrix elements) should be small . In addition, the prospects of determining the Majorana $`CP`$ phases increase with increasing solar neutrino mixing angle.
The quest for the limit on the absolute neutrino mass scale will witness attacks both via direct kinematical searches and cosmological measurements. The KATRIN experiment, currently under construction in Germany, is scheduled to start taking data in 2008 and is sensitive to neutrino masses down to 0.2 eV. Further cosmological probes will test the sum of neutrino masses down to 0.1 eV within this decade . Additional data sets and novel experiments can reduce this number by a factor of two .
Finally, the last piece of information needed to construct completely the neutrino mass matrix is the ordering of the neutrino states – the sign of $`\mathrm{\Delta }m_{31}^2`$. If the sign was positive we would have a normal mass ordering, while for a negative sign we would have an inverted mass ordering. Typical approaches to identify sgn$`(\mathrm{\Delta }m_{31}^2)`$ rely on using matter effects. Since the earth matter effect for modest baselines, and therefore modest matter densities, depends crucially on the size of the mixing angle $`\theta _{13}`$, it becomes exceedingly difficult to study the mass hierarchy as $`\theta _{13}`$ becomes small. If $`\theta _{13}`$ was large, close to its current limit, there could be a small chance of measuring the mass hierarchy using the synergies between T2K and NO$`\nu `$A . However, the measurement would still not be very unambiguous and one would need either a beta beam facility or a neutrino factory for the hierarchy determination . Very large matter effects in the 1–3 channel are expected for supernova neutrinos. Therefore, a supernova neutrino signal could in principle be used to determine the sign of $`\mathrm{\Delta }m_{31}^2`$ . Resonant matter effects in the 1–3 channel are also encountered by atmospheric neutrinos as they cross large baselines in their passage through the earth. This can be exploited to probe the neutrino hierarchy both in water Cerenkov and large magnetized iron calorimeter detectors . Recently some novel ways of probing the mass hierarchy requiring very precise measurements and using the “interference terms” between the different oscillation frequencies have been proposed (see also ). However, it is understood that among all the neutrino parameters, the mass hierarchy determination is expected – together with the determination of the Majorana phases – to be the most challenging for the future experiments. In this respect the neutrinoless double beta decay is of some interest, since the question of distinguishing the neutrino mass hierarchy can be answered by 0$`\nu \beta \beta `$. This issue will be discussed in the following Section.
## 3 Distinguishing the Neutrino Mass Schemes
In this Section we present the phenomenology of neutrinoless double beta decay in terms of the oscillation parameters $`\mathrm{\Delta }m_{31}^2`$, $`\mathrm{\Delta }m_{21}^2`$, $`\mathrm{sin}^2\theta _{12}`$ and $`\mathrm{sin}^2\theta _{13}`$. In particular we will look how easy or difficult it would be to distinguish the different neutrino mass schemes if we have a signal or a significantly improved limit for $`0\nu \beta \beta `$.
For the NH scheme, for $`m_1m_2m_3`$ and therefore assuming that $`m_1`$ can be neglected, we have
$`m^{\mathrm{NH}}\left|\sqrt{\mathrm{\Delta }m_{21}^2}\mathrm{sin}^2\theta _{12}\mathrm{cos}^2\theta _{13}+\sqrt{\mathrm{\Delta }m_{31}^2}\mathrm{sin}^2\theta _{13}e^{2i(\beta \alpha )}\right|.`$ (16)
The maximal value, $`m_{\mathrm{max}}^{\mathrm{NH}}`$, is achieved when both terms add up, which corresponds to $`\alpha \beta =\pm 2\pi ,\pm \pi `$ or 0. The minimal value of $`m`$ corresponds to $`\alpha \beta =\pm \pi /2,3\pi /2,5\pi /2`$ or $`7\pi /2`$, resulting in (partial) cancellation between the two terms in Eq. (16). For non–zero values of $`m_1`$ between $`10^2`$ and $`10^3`$ eV there can be complete cancellation resulting in a vanishing effective mass. This will be subject of Section 5.
For the IH scheme, assuming that $`m_3m_1<m_2`$ and neglecting $`m_3`$, we have
$`m^{\mathrm{IH}}\sqrt{|\mathrm{\Delta }m_{31}^2|}\mathrm{cos}^2\theta _{13}\sqrt{1\mathrm{sin}^22\theta _{12}\mathrm{sin}^2\alpha }.`$ (17)
Here the maximal (minimal) value is obtained when $`\alpha =0,\pi `$ or $`2\pi `$ ($`\alpha =\pi /2,3\pi /2,5\pi /2`$ or $`7\pi /2`$).
Finally, for the QD mass spectrum
$`m^{\mathrm{QD}}m_0\left|(\mathrm{cos}^2\theta _{12}+\mathrm{sin}^2\theta _{12}e^{2i\alpha })\mathrm{cos}^2\theta _{13}+e^{2i\beta }\mathrm{sin}^2\theta _{13}\right|.`$ (18)
If the terms proportional to $`e^{2i\alpha }`$ and $`e^{2i\beta }`$ add up (i.e., when $`\alpha `$ and $`\beta `$ take values of $`0,\pi `$ or $`2\pi `$) we have the maximal value of $`m`$, which is then just $`m_0`$. On the other hand, the minimal $`m`$ is achieved when $`\alpha `$ and $`\beta `$ take values of $`\pi /2,3\pi /2,5\pi /2`$ or $`7\pi /2`$. We will use this to set limits on $`m_0`$ in Section 4.
Hence, for a given mass scheme, the effective mass depends on $`\mathrm{\Delta }m_{31}^2`$, $`\mathrm{\Delta }m_{21}^2`$, $`\mathrm{sin}^2\theta _{12}`$, $`\mathrm{sin}^2\theta _{13}`$ and on the Majorana phases (to be precise, on one of them, or on a combination thereof). In case of QD the absolute neutrino mass scale $`m_0`$ plays a decisive role as well. In the case of NH the smallest neutrino mass could be extremely important in deciding the degree of cancellation between the different terms, as discussed above. The lack of knowledge of most of those parameters means that we can not give definite predictions for the value of $`m`$ for any of the three extreme neutrino mass schemes. We can however give a range of $`m`$ for NH, IH and QD, namely:
$`0.0<m<0.007\mathrm{eV}`$ $`,\mathrm{NH}`$ (19)
$`0.006\mathrm{eV}<m<0.065\mathrm{eV}`$ $`,\mathrm{IH}`$ (20)
$`0.41(0.07)\mathrm{eV}<m<2.3(0.4)\mathrm{eV}`$ $`,\mathrm{QD}`$ (21)
In order to obtain those values, we took the ranges of the oscillation parameters from Eqs. (10,11,12,14) and inserted them in the expressions for the maximal and minimal values for $`m`$ in the three extreme schemes under consideration. For the case of QD we used the limit from the direct kinematical search and in brackets $`m_0<0.4`$ eV. Note that with the $`3\sigma `$ values the NH and IH cases slightly overlap. However, these ranges are expected to reduce remarkably in the future from more precise measurements of the solar and atmospheric neutrino parameters and from either a signal for a non–zero value of (or from better upper limits on) $`\mathrm{sin}^2\theta _{13}`$. Nevertheless, already at the present stage the possibility of distinguishing the mass spectra from each other opens up. In particular, the case of deciding between NH and IH is the most interesting one. Distinguishing the QD spectrum from NH or IH will be at most a consistency check since the common neutrino mass of or above 0.2 eV will be probed by either direct searches or cosmology.
We present in Figure 2 the lines of constant maximal and minimal $`m`$ for the cases of NH and IH, respectively, in the $`\mathrm{sin}^2\theta _{12}\mathrm{sin}^2\theta _{13}`$ plane. We fixed the values of $`\mathrm{\Delta }m_{31}^2=0.002`$ eV<sup>2</sup> and $`\mathrm{\Delta }m_{21}^2=8\times 10^5`$ eV<sup>2</sup>. The dotted (red) lines show the lines of constant minimum value of $`m`$ for the IH scheme ($`m_{\mathrm{min}}^{\mathrm{IH}}`$) and the solid (blue) lines give the lines of constant maximum value of $`m`$ for the NH scheme ($`m_{\mathrm{max}}^{\mathrm{NH}}`$). The Figure illustrates the well–known fact that $`m_{\mathrm{min}}^{\mathrm{IH}}`$ has a strong $`\mathrm{sin}^2\theta _{12}`$ dependence, while $`m_{\mathrm{min}}^{\mathrm{NH}}`$ depends on both $`\mathrm{sin}^2\theta _{12}`$ and $`\mathrm{sin}^2\theta _{13}`$. Also shown by the dot–dashed (black) line on the Figure are the current allowed values of $`\mathrm{sin}^2\theta _{12}`$ and $`\mathrm{sin}^2\theta _{13}`$ at the 3$`\sigma `$ level, obtained from a 2 parameter fit of the global oscillation data <sup>4</sup><sup>4</sup>4The $`3\sigma `$ limit $`\mathrm{sin}^2\theta _{13}<0.044`$ given in the Introduction and Section 2 was obtained in a 1 parameter fit of the global neutrino data and is therefore different from the one in Figure 2.. As an example, one can see from the plot that if we know that the mass ordering is normal and assume that the smallest neutrino mass is negligible, values of the effective mass larger than roughly 0.0053 eV are incompatible with the currently allowed values of $`\mathrm{sin}^2\theta _{12}`$ and $`\mathrm{sin}^2\theta _{13}`$. The same would be the case when we know that the mass ordering is inverted, assume that the smallest neutrino mass is negligible and the effective mass is larger (or smaller) than 0.023 (0.008) eV.
We can see from Eqs. (16) and (17) that the maximal and minimal values of $`m`$ are crucially dependent on the values of the mass squared differences, especially on $`\mathrm{\Delta }m_{31}^2`$. Since we are talking about a future measurement of the $`0\nu \beta \beta `$ signal, we take the expected futuristic uncertainties mentioned in Section 2 on $`\mathrm{\Delta }m_{31}^2`$ and $`\mathrm{\Delta }m_{21}^2`$ to plot the iso–range of $`m_{\mathrm{min}}^{\mathrm{IH}}`$ and $`m_{\mathrm{max}}^{\mathrm{NH}}`$ in Figures 2 and 4. We can see clearly from the Figures that the uncertainty on the predicted value of $`m`$ due to the uncertainty in the value of $`\mathrm{\Delta }m_{31}^2`$ and $`\mathrm{\Delta }m_{21}^2`$ would become very small in the next ten years. Therefore, in what follows, we will keep $`\mathrm{\Delta }m_{31}^2`$ and $`\mathrm{\Delta }m_{21}^2`$ fixed at $`\mathrm{\Delta }m_{31}^2=0.002`$ eV<sup>2</sup> and $`\mathrm{\Delta }m_{21}^2=8\times 10^5`$ eV<sup>2</sup>.
Any future positive signal for $`0\nu \beta \beta `$ will be able to distinguish the IH scheme from the NH scheme if the measured value of $`m`$ would be such that it could be explained by the predicted $`m`$ for just one of the schemes and not the other. In other words, if the difference between the predicted values for $`m`$ for the IH scheme and the NH scheme is larger than the error in the measured value of $`m`$, then it would be possible (assuming that the smallest neutrino mass is negligible) to experimentally find the neutrino mass hierarchy from a measurement of $`0\nu \beta \beta `$. We therefore show the difference in the predicted value of $`m`$ for the two schemes,
$`\mathrm{\Delta }m=m_{\mathrm{min}}^{\mathrm{IH}}m_{\mathrm{max}}^{\mathrm{NH}},`$ (22)
in left–hand panel of Figure 4 as a function of $`\mathrm{sin}^2\theta _{12}`$ for 4 fixed values of $`\mathrm{sin}^2\theta _{13}`$ and in right–hand panel of Figure 4 as a function of $`\mathrm{sin}^2\theta _{13}`$ for 4 fixed values of $`\mathrm{sin}^2\theta _{12}`$. We see that $`\mathrm{\Delta }m`$ displays a strong dependence on $`\mathrm{sin}^2\theta _{12}`$, whereas the dependence on $`\mathrm{sin}^2\theta _{13}`$ is rather moderate. In fact, it holds that
$`\mathrm{\Delta }m={\displaystyle \frac{\sqrt{\mathrm{\Delta }m_\mathrm{A}^2}\mathrm{cos}^2\theta _{13}}{1+\mathrm{tan}^2\theta _{12}}}\left(1\mathrm{tan}^2\theta _{12}\left(1+R+\mathrm{tan}^2\theta _{13}\right)\mathrm{tan}^2\theta _{13}\right),`$ (23)
where $`R`$ is the ratio of the solar and atmospheric $`\mathrm{\Delta }m^2`$ defined in Eq. (15). Figure 4 and Eq. (23) demonstrate the well–known fact that distinguishing the normal and inverted hierarchy is easier for smaller values of $`\mathrm{sin}^2\theta _{12}`$. The values of $`\mathrm{\Delta }m`$, which are of the order 0.01 eV, represent a value of the maximal experimental uncertainty which an experiment should have in order to be able to distinguish NH from IH.
Up to here we neglected the complications which arise from the nuclear matrix elements (NME) involved in the calculation of the rate of 0$`\nu \beta \beta `$. There is neither a consensus in the literature on how large this uncertainty is, nor if there really is a large uncertainty at all . Let us therefore take a conservative point of view and analyze the situation as a function of the NME uncertainty. If there is indeed no problem with the NME, the statements given up to this point apply.
Given the possibility of large uncertainties on the values of the nuclear matrix elements, it is plausible to ask if one could possibly extract the true hierarchy of neutrino masses from a positive signal for $`0\nu \beta \beta `$ in the future. To that end, we look for the difference between $`m_{\mathrm{min}}^{\mathrm{IH}}`$ and $`m_{\mathrm{max}}^{\mathrm{NH}}`$ after including the uncertainties coming from the nuclear matrix elements. A way to perform such an analysis has been developed in , and here we follow a very similar approach: in order to parametrize the uncertainty coming from the nuclear matrix elements, we define the $`0\nu \beta \beta `$ decay rate $`\mathrm{\Gamma }`$ measured in any experiment such that
$`\sqrt{\mathrm{\Gamma }}=xm,`$ (24)
where $`x`$ contains all other factors involved in the $`0\nu \beta \beta `$ decay, including the NME. Thus for any measured $`\mathrm{\Gamma }`$ the measured value of $`m`$ would be
$`m_{\mathrm{measured}}=\sqrt{\mathrm{\Gamma }}/x.`$ (25)
If $`y`$ is the smallest possible value for $`x`$, and if we assume that the uncertainty on $`x`$ originates solely from the uncertainty on the NME and is a factor of $`z`$ (with $`z>1`$), then the range of $`x`$ is given by $`y<x<zy`$. Therefore, the uncertainty on the measured value of $`m`$ that we have due to the uncertainty on the NME is given by
$`\mathrm{\Delta }m_{\mathrm{measured}}={\displaystyle \frac{z1}{z+1}}.`$ (26)
In order to be able to experimentally distinguish IH from NH with the help of a $`0\nu \beta \beta `$ measurement, we must have $`m_{\mathrm{min}}^{\mathrm{IH}}>m_{\mathrm{max}}^{\mathrm{NH}}`$ after including the uncertainty from the NME. Thus we need the condition
$`m_{\mathrm{min}}^{\mathrm{IH}}zm_{\mathrm{max}}^{\mathrm{NH}}>0`$ (27)
to ascertain the neutrino mass hierarchy. In the left and right panel of Figure 6 we plot the difference $`m_{\mathrm{min}}^{\mathrm{IH}}zm_{\mathrm{max}}^{\mathrm{NH}}`$ predicted as a function of $`\mathrm{sin}^2\theta _{12}`$ and $`\mathrm{sin}^2\theta _{13}`$, respectively, including the uncertainty coming from the NME. We show the plots assuming a range of values for the NME uncertainty $`z`$. As a function of $`\mathrm{sin}^2\theta _{12}`$ the quantity $`m_{\mathrm{min}}^{\mathrm{IH}}zm_{\mathrm{max}}^{\mathrm{NH}}`$ shows a similar behavior as for no uncertainty in the NME. The dependence on $`\mathrm{sin}^2\theta _{13}`$ is rather weak for small $`z`$, but becomes larger with increasing $`z`$. From Figure 6 one can read of the maximal experimental uncertainty for a given $`z`$, $`\mathrm{sin}^2\theta _{13}`$ and $`\mathrm{sin}^2\theta _{12}`$ for which one could still be able to distinguish the mass hierarchy. For instance, taking $`\mathrm{sin}^2\theta _{13}`$ close to its current limit and assuming $`z=2`$ and $`\mathrm{sin}^2\theta _{12}=0.3`$, we have $`m_{\mathrm{min}}^{\mathrm{IH}}zm_{\mathrm{max}}^{\mathrm{NH}}0.01`$ eV and the experimental uncertainty should not be larger than this value. For zero $`\theta _{13}`$ this value becomes larger by roughly $`40\%`$.
Up to now we considered the case of the smallest mass $`m_{\mathrm{sm}}`$ being zero, which would correspond to the neutrino mass matrix having a zero determinant, clearly a rather limited and special case. For a given set of oscillation parameters and the nuclear matrix elements uncertainty, it would be possible to distinguish between the NH and IH in a $`0\nu \beta \beta `$ experiment only if $`m_{\mathrm{sm}}`$ is less than a certain value. It is interesting to examine up to which values of the lightest neutrino mass one can distinguish NH from IH. Toward this exercise, we plot in Figure 6 the difference between the minimal value of $`m`$ for IH and the maximal value of $`m`$ for NH, $`m_{\mathrm{min}}^{\mathrm{IH}}m_{\mathrm{max}}^{\mathrm{NH}}`$, as a function of the smallest neutrino mass $`m_{\mathrm{sm}}`$. The NME uncertainty is also included in the analysis. For each of the 4 panels we assumed a certain “true” value of $`\mathrm{sin}^2\theta _{12}`$ with an optimistic uncertainty of 6% . For $`\mathrm{sin}^2\theta _{13}`$ we assumed that no positive signal is observed in any of the forthcoming experiments in the next ten years. Thus we allowed it to vary within $`0`$ and $`0.0025`$ . We see from the Figure that for $`\mathrm{sin}^2\theta _{12}=0.3`$ for example, when $`z=1.5`$ ($`z=1`$, viz. no NME uncertainty) and $`\mathrm{\Delta }m=0.01`$ eV, we can distinguish NH from IH if $`m_{\mathrm{sm}}`$ is smaller than 0.002 (0.004) eV. Of course, the value of $`\mathrm{\Delta }m`$ obtained like this is again an upper limit on the experimental uncertainty which would still allow distinguishing NH from IH for a non–zero smallest mass. If $`\mathrm{sin}^2\theta _{12}=0.26`$ the situation looks more promising. For the same value of $`\mathrm{\Delta }m=0.01`$ eV we can distinguish NH from IH if $`m_{\mathrm{sm}}`$ is smaller than 0.005 (0.01) eV if $`z=1.5`$ ($`z=1`$). Other cases are easily read off Figure 6.
For a larger uncertainty in the value of $`\mathrm{sin}^2\theta _{12}`$ and/or for a larger true value of $`\mathrm{sin}^2\theta _{13}`$, the chances of distinguishing the IH from NH become worse and in addition only work for very small values of $`m_{sm}`$.
We see however that in principle values of $`m_{\mathrm{sm}}`$ up to $`0.01`$ eV are allowed in order to distinguish the normal from the inverted neutrino mass hierarchy.
## 4 Limit on the Neutrino Mass
Another interesting aspect of neutrinoless double beta decay is the possibility to set a limit on the absolute scale of the neutrino mass. The current upper value for the effective mass is somewhere between 0.3 and 1 eV, where this range of course has its origin in the NME uncertainty. The indicated values correspond to the QD mass spectrum, on which we wish to focus in this Section. As indicated in Section 3, for a given common mass scale $`m_0`$, the lowest possible effective mass can be written as
$`m_{\mathrm{min}}^{\mathrm{QD}}=m_0\left(|U_{e1}|^2|U_{e2}|^2|U_{e3}|^2\right)={\displaystyle \frac{1\mathrm{tan}^2\theta _{12}2|U_{e3}|^2}{1+\mathrm{tan}^2\theta _{12}}}.`$ (28)
Hence, having a limit on the effective mass at hand, we can translate it into a limit on the neutrino mass. We write the experimental limit as
$`m^{\mathrm{exp}}=zm_{\mathrm{min}}^{\mathrm{exp}},`$ (29)
where $`m_{\mathrm{min}}^{\mathrm{exp}}`$ is defined as the limit on $`m`$ obtained by using the largest available NME and $`z1`$ encodes again the NME uncertainty. Then, the limit on the neutrino mass reads
$`m_0zm_{\mathrm{min}}^{\mathrm{exp}}{\displaystyle \frac{1+\mathrm{tan}^2\theta _{12}}{1\mathrm{tan}^2\theta _{12}2|U_{e3}|^2}}zm_{\mathrm{min}}^{\mathrm{exp}}f(\theta _{12},\theta _{13}).`$ (30)
We have introduced here a function $`f(\theta _{12},\theta _{13})`$ in this expression, which separates the information available from neutrino oscillation experiments from the information coming from 0$`\nu \beta \beta `$ and $`m_0`$. We think that this formulation might be helpful in understanding how neutrinoless double beta decay and the absolute neutrino mass scale are related. Currently the uncertainty on $`f(\theta _{12},\theta _{13})`$ is around 50%, $`1.9<f(\theta _{12},\theta _{13})<5.6`$. It is expected to reduce to $``$ 21%($``$ 9) % at $`3\sigma `$ if a low energy $`pp`$ solar neutrino experiment (reactor experiment at the survival probability minimum) would be built. The uncertainty depends only little on the value of $`\theta _{13}`$. From the current limit on the effective mass, $`m0.35z`$ eV, with the accepted value of $`z3`$ for <sup>76</sup>Ge (see ), and our current knowledge of $`f(\theta _{12},\theta _{13})`$, we can set a limit on $`m_0`$ of 5.6 eV, clearly weaker than the limit from tritium experiments. Only for close to vanishing NME uncertainty (i.e., $`z=1`$) we can reach values of $`m_0<2`$ eV which are then comparable to ones from direct searches.
In Figure 7 we show the iso–contours of $`m_0`$ predicted for the QD mass spectrum (cf. Eq. (30)) in the $`\mathrm{sin}^2\theta _{12}\mathrm{sin}^2\theta _{13}`$ plane. For each of the iso–contours we have assumed an illustrative value of $`zm_{\mathrm{min}}^{\mathrm{exp}}=0.1`$ eV. Also shown is the current 3$`\sigma `$ allowed area of $`\mathrm{sin}^2\theta _{12}`$ and $`\mathrm{sin}^2\theta _{13}`$. From the Figure we can see that for a measured value of $`zm_{\mathrm{min}}^{\mathrm{exp}}=0.1`$ eV, the constraint on the common mass scale for a QD spectrum would be: $`0.2\mathrm{eV}<m_0<0.6`$ eV. For any other measured value of $`zm_{\mathrm{min}}^{\mathrm{exp}}`$ corresponding to a QD mass scheme, we can obtain the corresponding constraints on $`m_0`$ by scaling the above limit suitably, i.e., for $`zm_{\mathrm{min}}^{\mathrm{exp}}=0.2`$ eV we would have $`0.4\mathrm{eV}<m_0<1.2`$ eV. The limits on $`m_0`$ given above are of course with our existing knowledge about $`\theta _{12}`$ and $`\theta _{13}`$. With reduction of the range of allowed values of $`\theta _{12}`$ and $`\theta _{13}`$, especially $`\theta _{12}`$, the constraints on $`m_0`$ are expected to reduce substantially, as discussed above. For instance, if $`f(\theta _{12},\theta _{13})`$ was known with an uncertainty of $`20\%`$, say $`2.7<f(\theta _{12},\theta _{13})<4.0`$, then we would have for $`zm_{\mathrm{min}}^{\mathrm{exp}}=0.1(0.2)`$ eV that $`0.3\mathrm{eV}<m_0<0.4`$ eV ($`0.6\mathrm{eV}<m_0<0.8`$ eV). Of course, if we have no signal for 0$`\nu \beta \beta `$, but just an upper limit on $`zm_{\mathrm{min}}`$, we have no longer an allowed range on $`m_0`$, but an upper limit corresponding to the largest value in the range. From the examples given above, one can note that for the QD mass spectrum, a measurement or a better constraint on $`m`$ will lead to a stronger limit on the absolute neutrino mass scale compared to the current limit from direct kinematical searches.
## 5 Implications of a Vanishing effective Mass
It is well–known that not only a measurement but also a non–measurement of neutrinoless double beta decay has significant influence on our knowledge of the unknown neutrino parameters. In this Section we wish to discuss a rarely studied subject, namely the impact that a very small upper limit on $`m`$ would have<sup>5</sup><sup>5</sup>5For a related earlier analysis, see .. An exactly vanishing $`m`$ would correspond to a texture zero in the neutrino mass matrix in the charged lepton basis – certainly an interesting feature. However, to prove/observe an exactly vanishing effective mass is a formidable task.
In case of an inverted hierarchy we know that there is a lower limit of $`m`$ which is — when using current $`3\sigma `$ values of the oscillation parameters — given by roughly 0.006 eV, see Eq. (20). Let us assume that neutrinos are Majorana particles, have an inverted ordering and that the effective mass takes has an upper limit of order 0.01 eV. This is a situation which might arise if we know from some other independent experiment that sgn$`(\mathrm{\Delta }m_\mathrm{A}^2)=1`$ and the next generation 0$`\nu \beta \beta `$–experiments do not find a signal corresponding to a non–vanishing $`m`$, but give an upper limit on $`m`$ still above the theoretical limit $`m_{\mathrm{min}}^{\mathrm{IH}}`$. Then we can infer from Eq. (17) the values of $`\mathrm{sin}^2\theta _{12}`$ and the Majorana phase $`\alpha `$ which are still compatible with the data. In Fig. 9 we display the result for an experimental upper limit on $`m`$ of 0.04, 0.03, 0.02, 0.01 eV, taking $`m_3=0`$, $`\theta _{13}=0`$ and $`\mathrm{\Delta }m_{31}^2=0.002`$ eV<sup>2</sup>. We checked that the results are rather stable when we depart from these values within their current uncertainty. The current $`3\sigma `$ values of $`\mathrm{sin}^2\theta _{12}`$ are also indicated in the Figure. Allowed is the area to the right of the respective curves. For instance, if $`m<0.02`$ eV and $`\mathrm{sin}^2\theta _{12}=0.3`$ then $`\alpha `$ has to lie between 1.2 and 1.9, or $`\alpha \pi /2\pm 0.4`$. Alternatively, if $`m<0.02`$ eV and $`\mathrm{sin}^2\theta _{12}<0.28`$, then the IH case is ruled out.
Figure 9 shows a similar analysis for the case of NH with a smallest mass $`m_1=0`$. In our parametrization, as can be seen from Eq. (16), it is the combination $`\beta \alpha `$ which governs the destructive interference which leads to very small or zero $`m`$. We took very small upper limits on $`m`$ of 0.004, 0.003, 0.002 and 0.001 eV (hence this discussion will not become realistic within the next 10 years) and chose $`m_1=0`$, $`\mathrm{sin}^2\theta _{13}=0.04`$ $`\mathrm{\Delta }m_{21}^2=8\times 10^5`$ eV<sup>2</sup> and $`\mathrm{\Delta }m_{31}^2=0.002`$ eV<sup>2</sup>. Note that for very small values of $`\theta _{13}`$, corresponding to $`\mathrm{sin}^2\theta _{13}R\mathrm{sin}^2\theta _{12}0.01`$, the dependence on this combination of phases drops. In this part of the Figure it is the left of the respective curve which is allowed. For instance, for $`m<`$ 0.002 eV and a rather large value of $`\mathrm{sin}^2\theta _{12}=0.4`$ the phases have to be very close to $`\pi /2`$, namely $`\beta \alpha \pi /2\pm 0.2`$. Alternatively, if $`m<0.001`$ eV and $`\mathrm{sin}^2\theta _{12}>0.33`$, then the NH case is ruled out.
As mentioned earlier in Section 3 we note here that in the case of NH, $`m`$ could have a sizable dependence on the smallest neutrino mass state $`m_1`$. We plot therefore in Fig. 11 the same as in Fig. 9 but for a smallest neutrino mass of $`m_1=0.005`$ eV. We took again $`\mathrm{sin}^2\theta _{13}=0.04`$, $`\mathrm{\Delta }m_{21}^2=8\times 10^5`$ eV<sup>2</sup> and $`\mathrm{\Delta }m_{31}^2=0.002`$ eV<sup>2</sup> and chose the same four upper limits on $`m`$ as above. With a non–vanishing $`m_1`$ Eq. (16) is modified to
$`m\left|m_1\mathrm{cos}^2\theta _{12}\mathrm{cos}^2\theta _{13}+\sqrt{m_1^2+\mathrm{\Delta }m_{}^2}\mathrm{sin}^2\theta _{12}\mathrm{cos}^2\theta _{13}e^{2i\alpha }+\sqrt{\mathrm{\Delta }m_\mathrm{A}^2}\mathrm{sin}^2\theta _{13}e^{2i\beta }\right|`$ (31)
and thus it is no longer $`\beta \alpha `$ but both $`\alpha `$ and $`\beta `$ which play a role. In Figure 11 we fixed $`\alpha =\pi /2`$ which gives a negative sign for the second term in Eq. (31). We indicate the allowed (disallowed) values of the parameter regions in the Figure with A (D). If $`m<0.004`$ eV and $`m<0.003`$ eV the disallowed areas are already outside the currently allowed $`3\sigma `$ region of $`\mathrm{sin}^2\theta _{12}`$ for $`\alpha `$ fixed at $`\pi /2`$. To show the dependence of our results on $`\alpha `$, we plot in Figure 11 the allowed areas in the $`\alpha `$$`\beta `$ plane for four different fixed values of $`\mathrm{sin}^2\theta _{12}`$, keeping $`\mathrm{sin}^2\theta _{13}`$ fixed at 0.04 and for $`m_1=0.005`$ eV. Allowed are the areas within the respective curves.
An interesting aspect is when $`\theta _{13}`$ vanishes but $`m_1`$ is non–zero. Then the effective mass is a function of the phase $`\alpha `$ alone:
$`m\left|m_1\mathrm{cos}^2\theta _{12}+\sqrt{m_1^2+\mathrm{\Delta }m_{}^2}\mathrm{sin}^2\theta _{12}e^{2i\alpha }\right|.`$ (32)
In Figure 12 we show the areas in the $`\mathrm{sin}^2\theta _{12}`$$`\alpha `$ plane that would be allowed if we had an upper limit on $`m`$ of 0.001 eV (black line), 0.002 eV (red line), 0.003 eV (green line) and 0.004 eV (blue line). The areas outside the closed curves would be disallowed. We can see that for a limit of $`m`$=0.001 eV, all values of $`\alpha `$ could be possible while for $`m`$=0.004 eV, only the range $`\pi /2\pm 0.67`$ is allowed. This allowed range of $`\alpha `$ is not expected to improve with the reduction of the uncertainty on $`\mathrm{sin}^2\theta _{12}`$ since its $`\theta _{12}`$ dependence, as can be seen from the Figure, is very weak.
For values of the smallest mass of $`m_10.005`$ eV and $`\mathrm{sin}^2\theta _{12}0.3`$ both terms in Eq. (32) are of roughly the same magnitude (i.e., $`m_1m_2`$) and we could write
$$mm_{\mathrm{NH}}\sqrt{1\mathrm{sin}^22\theta _{12}\mathrm{sin}^2\alpha },\text{ where }m_{\mathrm{NH}}\sqrt{m_1^2+\mathrm{\Delta }m_{}^2}0.01\mathrm{eV}.$$
Note the similarity of this equation with Eq. (17) in case of IH.
To sum up, in certain cases it might be possible to significantly constrain the allowed values of the Majorana phases. Moreover, the allowed values of the solar neutrino oscillation parameter $`\theta _{12}`$ can be comparable to its current known range.
## 6 Neutrinoless Double Beta Decay, Future Neutrino Oscillation Data and the Identification of the Neutrino Mass Matrix
In this Section we wish to give a summary of typical neutrino mass matrices available in the vast literature. Predictions for and correlations between the neutrino observables are implied by each of the candidates and can be used to distinguish them. In particular, knowledge of the neutrino mass spectrum, $`R`$, $`U_{e3}`$, $`\theta _{23}`$ and of course the effective mass are then very helpful for identifying the correct mass matrix. A complete study of all possibilities is (and maybe can) not performed, however, we feel that the majority of the models will produce at the end of the day a mass matrix which will more or less correspond to one of the examples given here. For instance, there are many symmetries which eventually lead to a $`\mu \tau `$ symmetric mass matrix , see for instance Refs. for some examples. For more models and details we refer to the excellent reviews available . We should remark that for the given examples, unless otherwise stated, the charged lepton mass matrix $`m_{\mathrm{}}`$ is diagonal, i.e., for the matrix $`U_{\mathrm{}}`$ diagonalizing $`m_{\mathrm{}}m_{\mathrm{}}^{}`$ holds that $`U_{\mathrm{}}=\mathrm{𝟙}`$, and therefore the PMNS matrix is just $`U=U_{\mathrm{}}^{}U_\nu =U_\nu `$, where $`U_\nu `$ diagonalizes $`m_\nu `$. Further note that most correlations focus on the quantities $`m`$, $`R`$, $`\theta _{23}`$, $`|U_{e3}|`$ or $`\theta _{12}`$. Interplay with the parameters describing $`CP`$ violation is rarely studied.
Let us focus first on the normal hierarchy, for which there is a very successful and well–known texture, namely
$`m_\nu =m_0\left(\begin{array}{ccc}aϵ^2& bϵ& dϵ\\ & e& f\\ & & g\end{array}\right).`$ (36)
Here $`a,b,d,e,f,g`$ are complex parameters of order one and $`ϵ`$ is a real and small parameter typically of order of the Cabibbo angle $`\lambda 0.22`$. Such a mass matrix preserves at leading order the lepton charge $`L_e`$ and can be obtained already by simple models based on $`U(1)`$ charges or through sequential right–handed neutrino dominance . A correlation resulting from Eq. (36) is
$`m=c_1\sqrt{\mathrm{\Delta }m_\mathrm{A}^2}|U_{e3}|^2\text{ with }|U_{e3}|=c_2R,`$ (37)
where $`c_{1,2}`$ are functions of the order one parameters and naturally also of order one. With the parameters in the $`\mu \tau `$ block unspecified, the atmospheric neutrino mixing angle typically deviates sizable from its maximal value $`\pi /4`$, i.e., $`\theta _{23}=\pi /4c_3\sqrt{R}`$.
The discussion is also possible in the context of $`\mu \tau `$ symmetric mass matrices . Since their predictions include maximal $`\theta _{23}`$ and zero $`\theta _{13}`$, one is interested in breaking schemes of the symmetry. Breaking the symmetry in the $`e`$ sector of $`m_\nu `$ leads to the same correlation as in Eq. (37) but atmospheric neutrino mixing is much closer to $`\pi /4`$. On the other hand, breaking the $`\mu \tau `$ symmetry in the $`\mu \tau `$ sector of $`m_\nu `$ leads to
$`mc_1\sqrt{\mathrm{\Delta }m_\mathrm{A}^2}|U_{e3}|,`$ (38)
where $`|U_{e3}|`$ is now very small, $`|U_{e3}|=c_2R`$, and the deviation from maximal atmospheric neutrino mixing is again sizable, i.e., of order $`\sqrt{R}`$. Another interesting breaking scenario of $`\mu \tau `$ symmetric mass matrices is when the $`e\mu `$ element is not identical to, but the complex conjugate of the $`e\tau `$ element . Then one finds that $`U_{e3}0`$ but Re$`U_{e3}=0`$, i.e., there is maximal $`CP`$ violation with $`\delta =\pm \pi /2`$. A special case of $`\mu \tau `$ symmetric mass matrices, for which similar statements will apply, is given when $`\mathrm{tan}^2\theta _{12}=0.5`$, which then corresponds to the so–called tri–bimaximal mixing scheme . “Bimaximal” scenarios with their prediction $`\mathrm{tan}^2\theta _{12}=1`$ typically require contributions from the charged leptons, see below.
Perturbing a zeroth order mass matrix that leads to zero $`U_{e3}`$ and maximal $`\theta _{23}`$ with a random matrix containing small entries of the same order $`ϵ`$ , will lead to a semi–hierarchical mass spectrum (i.e., $`m_2>m_1`$ and not $`m_2m_1`$) with a somewhat larger $`m`$ and also to sizable deviations both from $`U_{e3}=0`$ and $`\theta _{23}=\pi /4`$ .
Other candidates studied frequently in the literature are minimal models in the sense of having zeros in $`m_\nu `$ or a minimal number of parameters in a see–saw context .
Table 2 summarizes our collection of the typical mass matrices for the normal mass hierarchy. Once we know that the neutrino mass spectrum is hierarchical, an information that could be provided by a negative search for 0$`\nu \beta \beta `$ in future experiments and knowing that sgn$`(\mathrm{\Delta }m_{31}^2)=+1`$, we can sort out the correct mass matrix when we have precise information on the oscillation parameters.
In case of the inverted hierarchy, summarized in Table 2, the most stable candidate “theory” corresponds to a mass matrix conserving the flavor charge $`L_eL_\mu L_\tau `$ . Such a matrix is given by
$`m_\nu =m_0\left(\begin{array}{ccc}0& \mathrm{cos}\theta & \mathrm{sin}\theta \\ & 0& 0\\ & & 0\end{array}\right),`$ (42)
predicting one massless neutrino, zero $`\theta _{13}`$, zero $`m`$, maximal solar neutrino mixing with $`\mathrm{\Delta }m_{}^2=0`$ and $`\theta _{23}=\theta `$. Barring extreme fine–tuning, it is impossible to perturb the structure of Eq. (42) such that the measured values of $`\theta _{12}`$ and $`\mathrm{\Delta }m_{}^2`$ are both accommodated in accordance with the data. Typically, after breaking it will hold that $`\mathrm{tan}^2\theta _{12}1\pm R/2`$ ($`\mathrm{sin}^2\theta _{12}\frac{1}{2}\pm R/4`$), whereas experimentally $`\mathrm{sin}^2\theta _{12}\frac{1}{2}\sqrt{R}`$ is required. An exception is a see–saw model in which the perturbations at high energy have the same order of magnitude as the terms allowed by the symmetry, see . In that example, $`U_{e3}`$ and the smallest mass state remain zero. Another way out is to take corrections from the charged lepton sector into account . The matrix $`U_{\mathrm{}}`$, which diagonalizes $`m_{\mathrm{}}m_{\mathrm{}}^{}`$, is multiplied from the left to $`U_\nu `$, which is the matrix diagonalizing $`m_\nu `$ from Eq. (42). The latter has to be perturbed in order to generate a non–zero $`\mathrm{\Delta }m_{}^2`$. Given the observation that the deviation from maximal solar neutrino mixing is determined by the Cabibbo angle, i.e., $`\pi /4\theta _{}\lambda 0.22\theta _C`$ , one assumes a CKM–like structure of $`U_{\mathrm{}}`$ <sup>6</sup><sup>6</sup>6Note that the so–called Quark–Lepton–Complementarity, i.e., the exact relation $`\theta _{12}+\theta _C=\pi /4`$ is a special case of this procedure.. The typical result of the described Ansatz lies in a correlation between the solar neutrino mixing and $`U_{e3}`$, namely $`\mathrm{sin}^2\theta _{12}\frac{1}{2}\mathrm{cos}\varphi \mathrm{cot}\theta _{23}|U_{e3}|`$ , where $`\varphi `$ is a phase appearing in $`m_\nu `$, which can be given by the Dirac phases as measurable in oscillation experiments. This depends however on the breaking in $`m_\nu `$. Nevertheless, both $`|U_{e3}|`$ and $`\mathrm{tan}^2\theta _{12}`$ are expected to be close to their current upper limits. Assuming minimal breaking in $`m_\nu `$ one can show that there are remarkable correlations between the observables, such as $`m\sqrt{\mathrm{\Delta }m_\mathrm{A}^2}\left|\mathrm{cos}2\theta _{12}+4iJ_{CP}/\mathrm{sin}^2\theta _{23}\right|`$, where $`J_{CP}`$ is the Jarlskog invariant for $`CP`$ violating effects in neutrino oscillations, which is proportional to $`|U_{e3}|/4`$. This represents one of the few examples where the $`CP`$ phases are part of the predicted correlations, depending however on the details of the model. The correlation $`\mathrm{tan}^2\theta _{12}14|U_{e3}|`$ is however independent of the breaking in $`m_\nu `$ and relies only on bi–large $`U_\nu `$ . Hence, the two discussed possibilities incorporating $`L_eL_\mu L_\tau `$ predict values for $`m`$ of similar size but predict either zero or large $`U_{e3}`$.
There is another zeroth order scheme for the inverted hierarchy, namely
$`m_\nu =m_0\left(\begin{array}{ccc}1& 0& 0\\ & 1/2& 1/2\\ & & 1/2\end{array}\right).`$ (46)
Zero $`U_{e3}`$ and maximal $`\theta _{23}`$ are predicted and the two leading mass states have the same sign. One can perturb this structure and the result is that both $`U_{e3}`$ and $`\theta _{23}\pi /4`$ are very small quantities, but the effective mass is larger than in case of matrices based on the conservation of $`L_eL_\mu L_\tau `$.
The dependence of $`m`$ on $`U_{e3}`$ goes from none to sizable to little in the three cases discussed, cf. with Table 2. So, if we know that the mass ordering is inverted, i.e., sgn$`(\mathrm{\Delta }m_{31}^2)=1`$ and the effective mass is of order $`\sqrt{\mathrm{\Delta }m_\mathrm{A}^2}`$, then we basically only have to distinguish between $`L_eL_\mu L_\tau `$ from Eq. (42) and the matrix given in Eq. (46). This can be done by knowing how to close $`m`$ is to $`\sqrt{\mathrm{\Delta }m_\mathrm{A}^2}`$. More precision data on $`U_{e3}`$ and $`\theta _{23}`$ will fix the details of the model.
Now we turn to quasi–degenerate neutrinos. Usually, in such scenarios both the normal and the inverted mass ordering can be accommodated. Typical examples are given in Table 3. There are matrices with two zero entries which are compatible with quasi–degenerate neutrinos and which predict an interesting interplay of observables . Anarchical mass matrices do in general not generate extreme mixing angles and therefore typically result in large deviations from zero $`U_{e3}`$ and maximal $`\theta _{23}`$, i.e., of order $`R`$.
Possibilities to obtain QD neutrinos are models based on flavor democracy , of which we quote in Table 3 one interesting recent and typical example . In such models both $`U_\nu `$ and $`U_{\mathrm{}}`$ are required in order to reproduce the neutrino data, resulting in a dependence of the neutrino observables on the charged lepton masses. This makes the predictions also a function of the democracy breaking scenario, therefore somewhat model–dependent. However, the smallness of $`|U_{e3}|`$ in such models can be attributed to the small ratios of the charged lepton masses.
One can also upgrade a hierarchical mass spectrum (e.g. resulting from sequential dominance) to a quasi–degenerate one by adding a term proportional to the unit matrix (e.g. connected to a $`SO(3)`$ symmetry) to it . The neutrino mixing angles display approximately the same behavior as in the NH case (i.e., as from Eq. (36)), however, the effective mass is close to the common mass scale, i.e., there is little cancellation in $`m`$. Moreover, the larger $`m_0`$ the smaller the phases .
Some attention has recently been caught by a model based on the discrete symmetry $`A_4`$ . Applying the most general SUSY threshold corrections to the resulting mass matrix leads to predictions such as very large atmospheric mixing, purely imaginary $`U_{e3}`$ with an absolute value of order $`\sqrt{R}`$ and $`m=m_0`$. Similarly, one can use a simple Abelian symmetry corresponding to the conservation of the flavor symmetry $`L_\mu L_\tau `$ and perturb the resulting mass matrix<sup>7</sup><sup>7</sup>7Note that this matrix is a special case of $`\mu \tau `$ symmetry.
$`m_\nu =m_0\left(\begin{array}{ccc}\mathrm{cos}\theta & 0& 0\\ & 0& \mathrm{sin}\theta \\ & & 0\end{array}\right),`$ (50)
to which the $`A_4`$ model corresponds when $`\theta =\pi /4`$. The deviations from zero $`U_{e3}`$ and maximal $`\theta _{23}`$ are small and inverse proportional to $`m_0^2`$. Note that interestingly most of the models presented here have $`mm_0`$, i.e., the mass parameters hopefully measured in direct laboratory searches and in neutrinoless double beta decay experiments should be almost identical.
Another zeroth order scheme for quasi–degenerate neutrinos corresponds to a matrix of the form
$`m_\nu =m_0\left(\begin{array}{ccc}0& 1/\sqrt{2}& 1/\sqrt{2}\\ & (1+\eta )/2& (1+\eta )/2\\ & & (1+\eta )/2\end{array}\right).`$ (54)
Exact bimaximal mixing is predicted and breaking of the matrix is also required to generate splittings between the mass states. Similar statements as for the matrix Eq. (42) hold, i.e., without extreme fine–tuning it is impossible to generate deviations from maximal solar neutrino mixing sizable enough not to be in conflict with the data. Hence, the contribution from the charged lepton sector are required, which will lead to similar correlations as for $`L_eL_\mu L_\tau `$ discussed above, such as $`\mathrm{tan}^2\theta _{12}14|U_{e3}|`$. The correlation between the other observables depends strongly on the breaking. Since however the $`ee`$ element is zero in Eq. (54), the effective mass is expected to be small, i.e., $`mm_0\mathrm{cos}2\theta _{12}`$.
Hence, the identification of the mass matrix in case of a QD spectrum can be performed when information also from direct kinematical or from cosmological measurements is available. With the expected future precision data the identification of the mass matrix should also be possible.
To sum up this Section, a complete determination of the all neutrino observables and consequent identification of correlations will be able to discriminate between the various possible models. At least some of the possible candidates will be ruled out. Note finally that the values of the two Majorana phases, whose determination turned out to be an extremely challenging task, is not really required in order to identify the mass matrix.
## 7 Conclusions
We have analyzed some aspects of the connection between neutrinoless double beta decay, the neutrino mass scale and neutrino oscillation parameters. In particular, we concentrated on the question whether the expected future precision on the oscillation parameters simplifies certain important consequences of a measurement of, or improved upper limit on, 0$`\nu \beta \beta `$. We first summarized the current situation of the determination of the nine physical parameters of the neutrino mass matrix and the prospects of improving our knowledge about them with future experiments. Then we analyzed how 0$`\nu \beta \beta `$ can help in distinguishing the normal mass ordering from the inverted one. We included the nuclear matrix element uncertainty in the analysis and pointed attention to the fact that distinguishing NH from IH depends on the value of the smallest neutrino mass $`m_{\mathrm{sm}}`$. We analyzed this point and found that in principle values of $`m_{\mathrm{sm}}=0.01`$ eV are allowed, where of course the issue of the nuclear matrix elements can significantly spoil this possibility.
Then we investigated inasmuch 0$`\nu \beta \beta `$ can be used to set a limit on the neutrino mass scale $`m_0`$ in case of a quasi–degenerate spectrum. We argued how the information from oscillation data and from 0$`\nu \beta \beta `$ can be separated. Current limits on $`m_0`$ are weaker than direct ones but can be significantly improved with future 0$`\nu \beta \beta `$–experiments and more precise knowledge of the oscillation parameters $`\mathrm{sin}^2\theta _{12}`$ and $`\mathrm{sin}^2\theta _{13}`$.
Next we studied the implications of a very small, or even vanishing effective mass. Still insisting that neutrinos are Majorana particles, one can put in this case interesting constraints on parameters, in particular on $`\mathrm{sin}^2\theta _{12}`$ and the Majorana phases.
Finally, we tried to perform a scan through the literature and identified neutrino mass matrices, which tend to arise frequently in many different models. We listed their predictions and correlations and argued how future precision date of the oscillation parameters and information from neutrinoless double beta decay can help to identify the true neutrino mass matrix or at least to sort out many unsuccessful ones.
Acknowledgments
This work was supported by the “Deutsche Forschungsgemeinschaft” in the “Sonderforschungsbereich 375 für Astroteilchenphysik” and under project number RO–2516/3–1 (W.R.) and PPARC grant number PPA/G/O/2002/00479 (S.C.). |
warning/0506/hep-th0506018.html | ar5iv | text | # Corrections to the Entropy in Higher Order Gravity
## I Introduction
During the late $`1960`$’s and early $`1970`$’s an intense research activity in the field of black hole physics lead to the discovery of the laws of black hole mechanics. Soon, it was realized a striking resemblance with the laws of thermodynamicsbeke ; hawking . In this context Bekenstein proposed that the entropy in black holes is proportional to the surface area of the event horizon. Moreover, Hawking found that black holes generate also a thermal radiation, due to quantum pair production in their gravitational potential gradient and the presence of the event horizon. The temperature of this radiation is given in terms of the surface gravity at the event horizon, $`T_H=\mathrm{}\kappa /2\pi `$. An immediate consequence of this identification of the temperature is that the proportionality constant between the entropy and the area is fixed, $`S_{BH}=A/4G`$.
Althought this knowledge about the black hole thermodynamics is well established, we do not have a complete answer to this question in the statistical mechanics framework. However, there has been a substantial progress in identifying the microscopic degrees of freedom responsible for the Bekenstein-Hawking entropy. This advance has come from string theoryvafa and loop quantum gravityashtekar . In the case of string theory it is present a massless spectrum including the graviton, and at low energy it gives supergravity effectives theories. Black Holes therefore appear as classical solutions of low energy string theory.
The next step in this line of research has been the study of the leading correction to entropy. A common characteristic of the different approaches is the proportionality to $`\mathrm{ln}S_{BH}`$ kaul ,carlip ,solodukin . However, the proportionality constant does not exhibit the same universality. It can be shown that logarithmic corrections to thermodynamic entropy arise in all thermodynamic systems when small stable fluctuations around equilibrium are taken into accountdas . The stability condition is equivalent to the specific heat being positive.
On the other hand the study of the thermodynamic properties of black holes has been extended to higher order gravity theoriesmyers . Within these theories there is a special class of gravitational actions, of higher order in the curvature, known as Lovelock gravity lovelock . Lovelock gravity is exceptional in the sense that although it contains higher powers of the curvature in the Lagrange density, the resulting equations of motion contain no more than second derivatives of the metric. It is also a covariant and ghost free theory as it happens in the case of Einstein’s General Relativity.
An important result that was found in the thermodynamic context is that the area law is a peculiarity of the Einstein-Hilbert theory scan . These facts motivates a deeper study of the thermodynamics of the black hole solutions of such exotic theories setare odintsov ,eaacb . In this paper we will study the corrections to the entropy for the black holes solutions of Lovelock gravity. We shall first briefly review such a formulation.
## II Higher Dimensional Gravity
The Lanczos-Lovelock action is a polynomial of degree $`[d/2]`$ in the curvature, which can be expressed in the language of forms as scan
$$I_G=\kappa \underset{m=0}{\overset{[d/2]}{}}\alpha _mL^{(m)},$$
(1)
where $`\alpha _m`$ are arbitrary constants, and $`L^{(m)}`$ is given by
$$L^{(m)}=ϵ_{a_1\mathrm{}a_d}R^{a_1a_2}R^{a_{2m1}a_{2m}}e^{a_{2m+1}}e^{a_d},$$
(2)
with $`R^{ab}`$ being the Riemann curvature two-forms given by
$$R^{ab}=d\omega ^{ab}+\omega _c^aw^{cb}.$$
(3)
Here $`w_{ab}`$ are the spin connection one-forms and $`e^a`$ the vielbein. A wedge product between forms is understood throughout.
The corresponding field equations can be obtained varying with respect to $`e^a`$ and $`w^{ab}`$. In scan the expression for the coefficients $`\alpha _m`$ was found requiring the existence of a unique cosmological constant. In such a case these theories are described by the action
$$I_k=\kappa \underset{p=0}{\overset{k}{}}c_p^kL^{(p)},$$
(4)
which corresponds to (1) with the choice
$$\alpha _p:=c_p^k=\{\begin{array}{cc}\frac{l^{2(pk)}}{(d2p)}\left(\begin{array}{c}k\\ p\end{array}\right)\hfill & ,pk\hfill \\ 0\hfill & ,p>k\hfill \end{array}$$
(5)
for the parameters, where $`1k[(d1)/2]`$. For a given dimension $`d`$, the coefficients $`c_m^k`$ give rise to a family of inequivalent theories, labeled by $`k`$ which represent the highest power of curvature in the Lagrangian. This set of theories possesses only two fundamental constants, $`\kappa `$ and $`l`$, related respectively to the gravitational constant $`G_k`$ and the cosmological constant $`\mathrm{\Lambda }`$ through
$$\kappa =\frac{1}{2(d2)\mathrm{\Omega }_{d2}G_k},$$
(6)
$$\mathrm{\Lambda }=\frac{(d1)(d2)}{2l^2}.$$
(7)
For black hole solutions that are asymptotically flat we consider the vanishing cosmological constant limit case. When $`l\mathrm{}`$ the only non-vanishing terms in Eq(4) is the kth one; therefore the action is obtained from Eq(1) with the choice of coefficients
$$\alpha _p:=\stackrel{~}{c}_p^k=\frac{1}{(d2k)}\delta _p^k,$$
(8)
in this case the action reads
$$\stackrel{~}{I}_k=\frac{\kappa }{(d2k)}ϵ_{a_1\mathrm{}a_d}R^{a_1a_2}R^{a_{2k1}a_{2k}}e^{a_{2k+1}}e^{a_d}.$$
(9)
Note that for $`k=1`$ the Einstein action without cosmological constant is recovered, while for $`k=2`$ we obtained the Gauss-Bonnet action,
$$I_2=\frac{(d2)!\kappa }{(d4)}d^dx\sqrt{g}(R_{\mu \nu \alpha \beta }R^{\mu \nu \alpha \beta }+4R_{\mu \nu }R^{\mu \nu }R^2).$$
(10)
Returning to the action (4) we remember that this set of theories possess asymptotically AdS black hole solutions given byscan ,
$`ds^2`$ $`=`$ $`\left(1+{\displaystyle \frac{r^2}{l^2}}\left({\displaystyle \frac{2G_kM+\delta _{d2k,1}}{r^{d2k1}}}\right)^{1/k}\right)dt^2+`$ (11)
$`{\displaystyle \frac{dr^2}{1+\frac{r^2}{l^2}\left(\frac{2G_kM+\delta _{d2k,1}}{r^{d2k1}}\right)^{1/k}}}+r^2d\mathrm{\Omega }_{d2}^2.`$
The black hole mass for any value of $`k`$ is a monotonically increasing function of the horizon radius $`r_+`$, which reads
$$M(r_+)=\frac{r_+^{d2k1}}{2G_k}\left(1+\frac{r_+^2}{l^2}\right)^k\frac{1}{2G_k}\delta _{d2k,1}.$$
(12)
The presence of the Kronecker delta within the metrics $`(\text{11})`$ signals the existence of two possible black hole vacua($`M=0`$) with different causal structures. The generic case, with $`d2k1`$, is
$`ds^2`$ $`=`$ $`\left(1+{\displaystyle \frac{r^2}{l^2}}\left({\displaystyle \frac{2G_kM}{r^{d2k1}}}\right)^{1/k}\right)dt^2+`$ (13)
$`{\displaystyle \frac{dr^2}{1+\frac{r^2}{l^2}\left(\frac{2G_kM}{r^{d2k1}}\right)^{1/k}}}+r^2d\mathrm{\Omega }_{d2}^2.`$
Analogously with the Schwarzschild-AdS metric, this set possesses a continuous mass spectrum, whose vacuum state is the AdS spacetime. The other case is obtained for odd dimensions, and it is a peculiarity of Chern-Simons theories. From $`(\text{11})`$ we obtain,
$`ds^2`$ $`=`$ $`\left(1+{\displaystyle \frac{r^2}{l^2}}\left(2G_kM+1\right)^{1/k}\right)dt^2+`$ (14)
$`{\displaystyle \frac{dr^2}{1+\frac{r^2}{l^2}\left(2G_kM+1\right)^{1/k}}}+r^2d\mathrm{\Omega }_{d2}^2.`$
Here, the black hole vacuum differs from AdS spacetime.
## III Canonical Formalism
In this section we reviewdas the derivation of the entropy in the canonical formalism. Looking for the entropy corrections it is considered the existence of small thermal fluctuations around the equilibrium. Then, we begin with the canonical partition function
$$Z(\beta )=_0^{\mathrm{}}\mathrm{\Omega }(E)e^{\beta E}𝑑E,$$
(15)
where $`\mathrm{\Omega }(E)`$ is the density of states, that can be obtained from the partition function doing an inverse Laplace transform
$$\mathrm{\Omega }(E)=\frac{1}{2\pi i}_{ci\mathrm{}}^{c+i\mathrm{}}Z(\beta )e^{\beta E}𝑑\beta =\frac{1}{2\pi i}_{ci\mathrm{}}^{c+i\mathrm{}}e^{S(\beta )}𝑑\beta ,$$
(16)
where
$$S(\beta )=\mathrm{ln}Z(\beta )+\beta E$$
(17)
is the entropy. The integral can be performed by the method of steepest descent around the saddle point $`\beta _0=1/T_0`$ such that $`S_0^{}=(S/\beta )_{\beta =\beta _0}=0`$. Here $`T_0`$ is the equilibrium temperature. Expanding the entropy around $`\beta _0`$, we have
$$S=S_0+\frac{1}{2}(\beta \beta _0)^2S_0^{\prime \prime }+\mathrm{}$$
(18)
Substituting (18) in (16) and integrating we obtain
$$\mathrm{\Omega }(E)=\frac{e^{S_0}}{\sqrt{2\pi S_0^{\prime \prime }}}.$$
(19)
Finally, using the Boltzmann’s formula, is obtained
$$𝒮=\mathrm{ln}\mathrm{\Omega }=S_0\frac{1}{2}\mathrm{ln}\left(S_0^{\prime \prime }\right)+\mathrm{}..$$
(20)
Here $`𝒮`$ is the entropy at equilibrium. This is to be distinguished from the function $`S(\beta )`$, which is the entropy at any temperature.
The logarithmic term can be transformed taking into account that $`S_0^{\prime \prime }`$ is the fluctuation of the mean squared energy, i.e,
$$S_0^{\prime \prime }=<E^2><E>^2,$$
(21)
and that the specific heat is $`C=(E/T)_{T_0}`$. Therefore
$$𝒮=S_0\frac{1}{2}\mathrm{ln}\left(CT^2\right).$$
(22)
This result apply to stable thermodynamic systems with small fluctuations around the equilibrium. The stability condition is equivalent to the specific heat being positive. On the other hand is assumed that the quantum fluctuations of the thermodynamics quantities under consideration are small. In other words, for black holes very close to extremality ($`T0`$), the fluctuation analysis ceases to be valid due to large quantum fluctuactionsdas .
## IV Corrections to the Entropy in Higher Order Gravity
### IV.1 Asymptotically AdS black hole solutions
Reviewing the thermodynamic properties of the black hole solution (13) we began with the expression for the Hawking temperature, that is
$$T_H=\frac{1}{4\pi k_Bk}\left((d1)\frac{r_+}{l^2}+\frac{d2k1}{r_+}\right),$$
(23)
where $`r_+`$ is the horizon radius. Note that for all $`k`$ such that $`d2k10`$ the temperature has the same behavior that the Schwarzschild-AdS black hole(k=1), that is: the temperature diverges at $`r_+=0`$. Also has a minimum at $`r_c`$ given by
$$r_c=l\sqrt{\frac{d2k1}{d1}},$$
(24)
and grows linearly for large $`r_+`$. Consequently we can calculate the specific heat $`C_k=\frac{M}{T}`$ as a function of $`r_+`$. Using (23) and (12) we obtain,
$$C_k=k\frac{2\pi k_B}{G_k}r_+^{d2k}\left(\frac{r_+^2+r_c^2}{r_+^2r_c^2}\right)\left(1+\frac{r_+^2}{l^2}\right)^{k1}.$$
(25)
Here the function $`C_k`$ has an unbounded discontinuity at $`r_+=r_c`$, signaling a phase transition. We will deal with black hole with horizon radius that satisfies the condition $`r_+>r_c`$, where the specific heat is positive and the correction formula (22) can be apply.
Finally we present the entropy function
$$S_k=k\frac{2\pi k_B}{G_k}_0^{r_+}r^{d2k1}\left(1+\frac{r^2}{l^2}\right)^{k1}𝑑r,$$
(26)
obtained from the Euclidean path integral formalismscan . Similar results are obtained in the Lagrangian formalism olea .
For simplicity we just perform the calculations for black holes with $`k=2`$. Therefore, for the entropy we get
$$S_2^{(0)}=\frac{4\pi k_B}{G_2}r_+^{d4}\left[\frac{1}{(d4)}+\frac{r_+^2}{(d2)l^2}\right]$$
(27)
where $`(0)`$ stands for the uncorrected entropy. In terms of this entropy the Hawking temperature and the specific heat are given by:
$`T_H`$ $`=`$ $`\left[{\displaystyle \frac{1}{8\pi k_Bl^2}}\left({\displaystyle \frac{4\pi k_B}{G_2l^2}}\right)^{1/(d2)}\right]\times `$
$`\times {\displaystyle \frac{\left[(d1)+\frac{(d5)l^2}{r_+^2}\right]}{\left[\frac{l^2}{(d4)r_+^2}+\frac{1}{d2}\right]^{1/(d2)}}}(S_2^{(0)})^{1/(d2)},`$
$$C_2=\left(\frac{r_+^2+r_c^2}{r_+^2r_c^2}\right)\left(1+\frac{l^2}{r_+^2}\right)\frac{1}{\frac{l^2}{(d4)r_+^2}+\frac{1}{(d2)}}S_2^{(0)}.$$
(29)
In the limit $`\mathrm{r}_+l`$, $`C_2`$ approaches the value
$$C_2=(d2)S_2^{(0)},$$
(30)
and for the entropy we get
$$𝒮_2=S_2^{(0)}\frac{d}{2(d2)}\mathrm{ln}S_2^{(0)}+\mathrm{}.$$
(31)
This result is identical to the one obtained for the AdS-Schwarzschild black holesdas .
We can also study the correction to the entropy near the transition point $`r_c`$ (after the minimum of $`C_2`$). Assuming $`r_+^2r_c^2l^2`$(small $`l`$), the entropy behaves as
$$S_2^{(0)}=\frac{4\pi k_B}{G_2}r_+^{d4}\left[\frac{1}{(d4)}+\frac{2(d3)}{(d2)(d1)}\right].$$
(32)
Consequently the relations between the thermodynamic quantities are
$$T_H=\frac{B}{A^{1/(d4)}}(S_2^{(0)})^{1/(d4)},$$
(33)
$$C_2=\frac{C}{A}S_2^{(0)},$$
(34)
where
$$A=\frac{4\pi k_B}{G_2}\left[\frac{1}{(d4)}+\frac{2(d3)}{(d2)(d1)}\right],$$
(35)
$$B=\frac{1}{8\pi k_B}\frac{1}{l^2}\left[(d1)+\frac{(d5)(d1)}{2(d3)}\right],$$
(36)
$$C=\frac{4\pi k_B}{G_k}\left[1+\frac{2(d3)}{d1}\right]\left[\frac{2(d3)}{(d1)}+\frac{(d5)}{(d1)}\right].$$
(37)
Therefore
$$𝒮_2=S_2^{(0)}\frac{d2}{2(d4)}\mathrm{ln}S_2^{(0)}+\mathrm{}.$$
(38)
It is interesting to note that the entropy correction in this limit is greater than the entropy correction found in (31).
### IV.2 Chern-Simons Black Holes
Now let us consider the Chern-Simons black holes with metric (14). The thermodynamic quantities can be obtained considering $`d2k1=0`$. So,
$$T_H=\frac{1}{4\pi k_Bk}\left((d1)\frac{r_+}{l^2}\right),$$
(39)
$$C_{CS}=k\frac{2\pi k_B}{G_k}r_+\left(1+\frac{r_+^2}{l^2}\right)^{k1},$$
(40)
and
$$S_k=k\frac{2\pi k_B}{G_k}_0^{r_+}\left(1+\frac{r^2}{l^2}\right)^{k1}𝑑r.$$
(41)
In this case the temperature is not divergent and the specific heat is a continuous monotonically increasing positive function of $`r_+`$
In the limit $`r_+l`$ we get
$$C_{CS}=(2k1)S_k^0$$
(42)
and
$$𝒮_k=S_k^{(0)}\frac{2k+1}{2(2k1)}\mathrm{ln}S_k^{(0)}+\mathrm{}..$$
(43)
Note that in this limit the results, as was expected, are similar to those found in the previous section.
On the other hand, for $`r_+l`$, we obtain
$$𝒮_k=S_k^{(0)}\frac{3}{2}\mathrm{ln}S_k^{(0)}+\mathrm{}..$$
(44)
Similar to the proportionality constant found in das for the BTZ black holebanados .
## V Conclusions.
In this paper we have calculated the entropy corrections for different kinds of black holes in the context of higher order gravity. The results are similar to those found in the literature despite the fact that the area law is not satisfied. Also is confirmed the lack of universality of the logarithmic prefactor. In the case of Chern-Simons black holes with small horizon radius we have found that the logarithmic prefactor does not depended of the dimension.
ACKNOWLEDGMENT:
The author thanks Prof. Sandro Silva e Costa at the Physics Department-UFMT for hospitality. This work was supported by Fundação de Amparo à Pesquisa do Estado de Mato Grosso (FAPEMAT) and Conselho Nacional de Desenvolvimento Científico e Tecnológico (CNPq). The authors also thank the organizers of the Conference ”$`100`$ years of Relativity” where these results were presented. |
warning/0506/hep-ph0506063.html | ar5iv | text | # 1 Diagram introducing corrections of order 𝑀² to 𝑚_ℎ².
## Abstract
Little Higgs models are formulated as effective theories with a cut-off of up to 100 times the electroweak scale. Neutrino masses are then a puzzle, since the usual see-saw mechanism involves a much higher scale that would introduce quadratic corrections to the Higgs mass parameter. We propose a model that can naturally accommodate the observed neutrino masses and mixings in Little Higgs scenarios. Our framework does not involve any large scale or suppressed Yukawa couplings, and it implies the presence of three extra (Dirac) neutrinos at the TeV scale. The masses of the light neutrinos are induced radiatively, they are proportional to small ($``$ keV) mass parameters that break lepton number and are suppressed by the Little Higgs cut-off.
UG-FT-187/05
CAFPE-57/05
A Little Higgs model of neutrino masses
F. del Águila, M. Masip and J.L. Padilla
CAFPE and Departamento de Física Teórica y del Cosmos
Universidad de Granada, E-18071, Granada, Spain
faguila@ugr.es, masip@ugr.es, jluispt@ugr.es
Introduction. The stability of the electroweak (EW) scale at the loop level has been the main motivation to search for physics beyond the standard model (SM) during the past 30 years. Namely, SM loops introduce quadratic corrections to the EW scale. If this scale is natural, consistent with the dynamics and not the result of an accidental cancellation between higher scales, then new physics must compensate the SM quadratic contributions. In particular, top quark corrections to the Higgs mass squared become of order $`(500`$ GeV$`)^2`$ for a cutoff $`2`$ TeV, what would suggest new contributions (from supersymmetry, extra dimensions,…) of the same order at the reach of the LHC.
Little Higgs (LH) ideas provide another very interesting framework with a scalar sector free of unsuppressed loop corrections. New symmetries protect the EW scale and define consistent models with a cutoff as high as $`\mathrm{\Lambda }10`$ TeV. Therefore, these models could describe all the physics to be explored in the next generation of accelerators. More precisely, in LH models the scalar sector has a (tree-level) global symmetry G that is broken spontaneously at a scale $`f1`$ TeV. The SM Higgses are then Goldstone bosons (GBs) of the broken symmetry, and remain massless and with a flat potential at that scale. Yukawa and gauge interactions break explicitly the global symmetry. However, the models are built in such a way that the loop diagrams giving non-symmetric contributions must contain at least two different couplings. This collective breaking keeps the Higgs sector free of quadratic top-quark and gauge contributions. At the same time, loops (and/or explicit non-symmetric terms in the scalar potential) give mass and quartic couplings to the Higgs, both necessary to break the SM symmetry (see for a recent review).
The inclusion of neutrino masses in this framework looks problematic. In the SM these masses require a new scale much larger than the EW one,
$$_{eff}=\frac{1}{2\mathrm{\Lambda }_\nu }h^{}h^{}LL+\mathrm{h}.\mathrm{c}.,$$
(1)
where $`h=(h^0h^{})`$ and $`L=(\nu e)`$ are, respectively, the SM Higgs and lepton doublets. At the EW phase transition $`h^0=v`$ and the neutrinos get their mass $`m_\nu =v^2/\mathrm{\Lambda }_\nu 0.1`$ eV, what implies $`\mathrm{\Lambda }_\nu 10^{14}`$ GeV. This effective (low-energy) scenario is simply realised using the see-saw mechanism . A SM singlet per family, $`n^c`$, is introduced with a large Majorana mass $`M`$ and sizeable Yukawa couplings $`\lambda _\nu `$ with the lepton doublets:
$$_\nu =\lambda _\nu h^{}Ln^c+\frac{1}{2}Mn^cn^c+\mathrm{h}.\mathrm{c}.$$
(2)
(the fields denote two-component left-handed spinors and family indices are omitted). The spectrum is then three (Majorana) fields of mass $`M`$ plus the observed low-energy neutrinos with mass $`m_\nu \lambda _\nu ^2v^2/M`$.
This scenario looks inconsistent in LH models, since the EW scale is not stable at the quantum level.
In particular, the diagram in Fig. 1 gives a large contribution to the Higgs mass proportional to $`M^2`$:
$$\mathrm{\Delta }m_h^2\frac{\lambda _\nu ^2}{8\pi ^2}\left(C_{UV}+\frac{1}{2}M^2\left(12\mathrm{log}M^2\right)\right),$$
(3)
where we have used dimensional regularization and $`C_{UV}`$ contains the ultraviolet divergence and renormalization-scale dependence. In principle LH models would avoid this type of corrections from heavy fields just because they are supposed to be effective theories valid only below a cutoff $`\mathrm{\Lambda }(4\pi )^2v`$. However, if the see-saw scale $`M`$ is at (or below) the cut-off $`\mathrm{\Lambda }`$, neutrino masses would be unacceptably large. Therefore, the problem would be how to generate the operator in Eq. (1) without introducing also a term $`\mathrm{\Delta }m_h^2h^{}h`$ of order $`M^2`$ in the low-energy Lagrangian. The option of not using the see-saw mechanism and assume that neutrino masses are not different in origin from the masses of quarks and leptons requires $`M0`$ and $`\lambda _\nu 10^{12}`$, a number that would demand an explanation (see below). Other scenarios for neutrino masses which do not involve a mass $`M1`$ TeV could be consistent with LH ideas. For example, the model in has $`M1`$ TeV and extra scalars that give masses to neutrinos but not to quarks and charged leptons. The symmetries and structure of this model (the scalars get VEVs of order MeV), however, seem difficult to accommodate in a simple LH model. Here we propose a LH framework for neutrino masses that does not require large masses nor suppressed Yukawa couplings. It involves a quasi-Dirac field per family at the TeV scale and small ($``$ keV) Majorana mass terms breaking lepton number. The SM neutrinos are then Majorana fields that get their masses and mixings at the loop level.
Cancellations in Little Higgs models. LH models are able to explain why the top quark does not introduce one-loop quadratic corrections to $`m_h^2`$. In all the cases this is achieved extending the global symmetries to the quark sector, so that the third generation appears in multiplets of $`SU(3)`$. In particular, the doublet $`Q=(tb)`$ becomes a triplet $`\mathrm{\Psi }_Q`$. In the same way, in order to build a consistent model of neutrino masses the global symmetry will be extended and the lepton doublets will become triplets.
Let us focus on the simplest LH model , although our arguments would be analogous in the original littlest Higgs model or other more complicated scenarios . Here the scalar sector contains two triplets, $`\varphi _1`$ and $`\varphi _2`$, of a global $`SU(3)_1\times SU(3)_2`$ symmetry:
$`\varphi _1`$ $``$ $`e^{i\theta _1^aT^a}\varphi _1,`$
$`\varphi _2`$ $``$ $`e^{i\theta _2^aT^a}\varphi _2,`$ (4)
where $`T^a`$ are the generators of $`SU(3)`$. To get gauge interactions, the diagonal combination of the two $`SU(3)`$ is made local:
$$\varphi _{1(2)}e^{i\theta ^a(x)T^a}\varphi _{1(2)}.$$
(5)
At the scale $`f`$ the scalar triplets get vacuum expectation values (VEVs) and break the global symmetry to $`SU(2)_1\times SU(2)_2`$. For simplicity, it is usually assumed identical VEVs for both triplets
$$\varphi _1=\left(\begin{array}{c}0\\ 0\\ f\end{array}\right),\varphi _2=\left(\begin{array}{c}0\\ 0\\ f\end{array}\right).$$
(6)
The initial 12 scalar degrees of freedom in the two triplets contain 10 GBs plus two massive fields. However, these VEVs also break the gauge symmetry $`SU(3)\times U(1)_\chi `$ to $`SU(2)_L\times U(1)_Y`$, a process that will absorb 5 of the GBs. All this becomes apparent if the two triplets are parameterized
$`\varphi _{1(2)}=`$ $`\mathrm{exp}\left\{{\displaystyle \frac{i}{f}}\left(\begin{array}{cc}& h^{}\\ h^{}& \eta ^{}\end{array}\right)\right\}\times `$ (14)
$`\mathrm{exp}\left\{+(){\displaystyle \frac{i}{f}}\left(\begin{array}{cc}& h\\ h^{}& \eta \end{array}\right)\right\}\left(\begin{array}{c}0\\ f+\frac{r_{1(2)}}{\sqrt{2}}\end{array}\right),`$
where $`h^{}`$ and $`h`$ are (complex) $`SU(2)`$ doublets and $`\eta ^{}`$, $`\eta `$, $`r_1`$ and $`r_2`$ are (real) $`SU(2)`$ singlets. At the scale $`f`$ $`h^{}`$ and $`\eta ^{}`$ are eaten by the massive vector bosons of $`SU(3)`$, $`r_{1,2}`$ get massive, and $`h`$ (and possibly $`\eta `$) are the SM Higgses. Therefore
$`\varphi _{1(2)}`$ $``$ $`\left(\begin{array}{c}0\\ f+\frac{r_{1(2)}}{\sqrt{2}}\end{array}\right)+()i(1+{\displaystyle \frac{r_{1(2)}}{f\sqrt{2}}})\left(\begin{array}{c}h\\ \eta \end{array}\right)`$ (22)
$`{\displaystyle \frac{1}{2}}(1+{\displaystyle \frac{r_{1(2)}}{f\sqrt{2}}})\left(\begin{array}{c}\eta h\\ h^{}h+\eta ^2\end{array}\right).`$
In this model the top-quark Yukawa sector includes a triplet $`\mathrm{\Psi }_Q=(QT)`$ and two singlets ($`t_1^c`$, $`t_2^c`$), and it is described by the Lagrangian
$`_t`$ $`=`$ $`\lambda _1\varphi _1^{}\mathrm{\Psi }_Qt_1^c+\lambda _2\varphi _2^{}\mathrm{\Psi }_Qt_2^c+\mathrm{h}.\mathrm{c}.`$ (23)
$``$ $`\lambda _t(h^{}Qt^c+fTT^c{\displaystyle \frac{1}{2f}}h^{}hTT^c)+\mathrm{h}.\mathrm{c}.,`$
where $`t^c=(i/\sqrt{2})(t_2^ct_1^c)`$, $`T^c=(1/\sqrt{2})(t_2^c+t_1^c)`$, and we have taken $`\lambda _1=\lambda _2=\lambda _t/\sqrt{2}`$. The one-loop quadratic corrections in Fig. 2 cancel, which reflects that if one of the $`\lambda _{1,2}`$ couplings is zero the global $`SU(3)_1\times SU(3)_2`$ symmetry would be exact in this sector.
Analogously, the model includes one lepton triplet $`\mathrm{\Psi }_L=(iLN)`$ (we use the $`i`$ phase to simplify the couplings) and one singlet $`n^c`$ per generation . The Lagrangian
$`_\nu `$ $`=`$ $`\lambda _\nu \varphi _1^{}\mathrm{\Psi }_Ln^c+\mathrm{h}.\mathrm{c}.`$ (24)
$``$ $`\lambda _\nu (h^{}Ln^c+fNn^c{\displaystyle \frac{1}{2f}}h^{}hNn^c)+\mathrm{h}.\mathrm{c}.`$
respects the global symmetries and does not generate one-loop quadratic corrections to $`m_h^2`$. When the Higgs $`h`$ gets a VEV $`v=174`$ GeV and breaks the EW symmetry, the Yukawa couplings induce mass terms and define the matrix
$$\begin{array}{cc}\begin{array}{ccc}\nu & N& n^c\end{array}& \\ \left(\begin{array}{ccc}0& \mathrm{\hspace{0.33em}\hspace{0.33em}0}& \lambda _\nu v\\ 0& \mathrm{\hspace{0.33em}\hspace{0.33em}0}& \lambda _\nu f\\ \lambda _\nu v& \lambda _\nu f& \mathrm{\hspace{0.33em}\hspace{0.33em}0}\end{array}\right)& \begin{array}{c}\nu \\ N\\ n^c\end{array}\end{array}.$$
(25)
The neutrino sector will then contain three Dirac fields of mass $`\lambda _\nu f`$ plus three massless neutrinos $`\nu ^{}\nu +v/fN`$. Several observations are here in order.
(i) The SM neutrinos are exactly massless at this stage. To give them masses one would need extra singlets (which could combine with them to define Dirac neutrinos) and/or new terms that break lepton number (which could introduce Majorana masses for the SM neutrinos).
(ii) The neutrino sector in Eq. (24) does not break the $`SU(3)_1\times SU(3)_2`$ symmetry: both $`\varphi _1`$ and $`\mathrm{\Psi }_L`$ are triplets of $`SU(3)_1`$ and singlets of $`SU(3)_2`$. At the loop level $`\lambda _\nu `$ will contribute to the (invariant) operator $`\varphi _1^{}\varphi _1`$ in the scalar potential, but it will not introduce quadratic corrections to $`m_h^2`$ (which would require a symmetry-breaking operator $`\varphi _1^{}\varphi _2`$). Quadratic corrections to $`m_h^2`$ would appear at one loop if we add to the Lagrangian a term $`\lambda _\nu ^{}\varphi _2^{}\mathrm{\Psi }_Ln^c`$, and at higher order if we combine $`\lambda _\nu `$ with other Yukawa and gauge couplings.
(iii) Any realistic LH model needs a mass term $`m_h^2\varphi _1^{}\varphi _2`$ and a quartic coupling $`\lambda (\varphi _1^{}\varphi _2)^2`$ to trigger EW symmetry breaking. These terms may appear at the loop level, from scalar couplings with the global symmetry-breaking sectors of the model . Once these terms are included, higher order corrections will induce the terms $`\lambda _\nu ^{}\varphi _2^{}\mathrm{\Psi }_Ln^c`$ in the Lagrangian.
A mechanism for neutrino masses in LH models. To obtain massive SM neutrinos we must then introduce additional singlets or break lepton number. It is easy to see that the first possibility requires very suppressed Yukawa couplings. An extra singlet (per generation) would make this sector identical to the top-quark sector in Eq. (23), with two Dirac fields of masses proportional to $`f`$ and $`v`$. The SM neutrinos would then be too heavy unless the corresponding Yukawa couplings are very suppressed. This escenario could be naturally realized in models with extra dimensions or in holographic models , where the Higgs appears as a composite particle of some strongly coupled dynamics.
Therefore, the scenario that we propose requires just one singlet $`n^c`$ per family and the breaking of lepton number. The simplest way to parameterize this breaking is through a Majorana mass for $`n^c`$, so we add the term $`\frac{1}{2}Mn^cn^c+\mathrm{h}.\mathrm{c}.`$ in the Lagrangian in Eq. (24). At lowest order the neutrino mass matrix would read
$$\begin{array}{cc}\begin{array}{ccc}\nu & N& n^c\end{array}& \\ \left(\begin{array}{ccc}0& \mathrm{\hspace{0.33em}\hspace{0.33em}0}& \lambda _\nu v\\ 0& \mathrm{\hspace{0.33em}\hspace{0.33em}0}& \lambda _\nu f\\ \lambda _\nu v& \lambda _\nu f& M\end{array}\right)& \begin{array}{c}\nu \\ N\\ n^c\end{array}\end{array}.$$
(26)
This matrix implies two massive states and one massless neutrino per generation. The massless field, however, will get a mass at the loop level. The diagram if Fig. 3 generates terms like
$`_1`$ $`=`$ $`{\displaystyle \frac{1}{2\mathrm{\Lambda }_\nu }}(\varphi _2^{}\mathrm{\Psi }_L)(\varphi _2^{}\mathrm{\Psi }_L)+\mathrm{h}.\mathrm{c}.`$ (27)
$``$ $`{\displaystyle \frac{1}{2\mathrm{\Lambda }_\nu }}(h^0\nu +fN)(h^0\nu +fN)+\mathrm{h}.\mathrm{c}.`$
that will induce masses for the SM neutrinos. If $`M<f`$ and $`m_h^2<(\lambda _\nu f)^2`$ and assuming diagonal couplings we obtain
$$\frac{1}{\mathrm{\Lambda }_\nu }\frac{\lambda M}{16\pi ^2f^2}\frac{x1\mathrm{log}x}{(x1)^2},$$
(28)
where $`xm_h^2/(\lambda _\nu f)^2`$ and $`\lambda `$ is the Higgs quartic coupling. Notice that $`1/\mathrm{\Lambda }_\nu `$ vanishes if $`M=0`$, since the term in Eq. (27) breaks lepton number. Combined with the terms in Eq. (24) it also breaks the global symmetries (it is proportional to the Higgs quartic coupling $`\lambda `$). Although the mechanism to generate $`\lambda `$ may be different in each LH model, $`\lambda `$ must be $`2m_h^2/v^2`$.
The complete neutrino mass matrix is then
$$\begin{array}{cc}\begin{array}{ccc}\nu & N& n^c\end{array}& \\ \left(\begin{array}{ccc}v^2/\mathrm{\Lambda }_\nu & vf/\mathrm{\Lambda }_\nu & \lambda _\nu v\\ vf/\mathrm{\Lambda }_\nu & f^2/\mathrm{\Lambda }_\nu & \lambda _\nu f\\ \lambda _\nu v& \lambda _\nu f& M\end{array}\right)& \begin{array}{c}\nu \\ N\\ n^c\end{array}\end{array}.$$
(29)
Its diagonalization gives two heavy neutrinos, $`n^c`$ and $`N^{}Nv/f\nu `$ (they define a quasi Dirac field), of mass $`\lambda _\nu f`$ plus a SM neutrino $`\nu ^{}\nu +v/fN`$ of mass (in the limit $`m_h^2(\lambda _\nu f)^2`$)
$$m_\nu 4v^2/\mathrm{\Lambda }_\nu v^2\frac{\lambda M}{4\pi ^2f^2}\mathrm{log}\frac{(\lambda _\nu f)^2}{m_h^2}$$
(30)
per family. To obtain $`m_\nu 0.1`$ eV the Majorana mass $`M`$ must be $`0.1`$ keV.
A first interesting consequence of this result (that distinguishes our framework from other scenarios for neutrino masses) is that the light masses tend to depend only logarithmically on the Yukawa couplings. Any difference or hierarchy in the Yukawa sector will appear softened by the logarithm in the neutrino masses.
In our model the large mixings in the Maki-Nakagawa-Sakata (MNS) matrix are obtained once the couplings $`\lambda _\nu `$ and the lepton number violating masses $`M`$ are allowed to be arbitrary $`3\times 3`$ matrices. An acceptable rate for FCNC processes (mediated at one loop by the TeV singlets) will require the alignment of the neutrino and the charged-lepton Yukawa couplings. If the charged-lepton mass matrix comes from terms $`(\lambda _e/\mathrm{\Lambda })\varphi _1\varphi _2\mathrm{\Psi }_Le^c`$ , the rotations diagonalizing $`\lambda _\nu `$ and $`\lambda _e`$ must coincide to a large extent. Then, the MNS matrix results from the diagonalization of Eq. (30) (generalized to non-diagonal couplings) in the basis where $`\lambda _\nu `$ is diagonal. Obviously, the allowed misalignment increases for larger values of $`f`$ (a detailed quantitative analysis will be presented elsewhere ).
The model that we propose provides an explicit realization of TeV-scale non-decoupled neutrinos, which may be observable at a large ($`e^+e^{}`$) collider . The fields $`N`$ have a mixing $`v/f`$ with the SM neutrinos that must be smaller than 0.07 to be consistent with current data . For $`N`$ masses of order TeV, single heavy neutrino production $`e^+e^{}N\nu eW\nu `$ could give a signal at CLIC if $`v/f0.005`$ .
Summary and discussion. In LH models the EW scale is protected from large quadratic corrections only at the one-loop level. Therefore, the framework can not naturally accommodate physics at scales larger than 10 TeV. Neutrino masses are then a puzzle, because in order to explain their size the effective low energy model must involve a much larger scale or very suppressed Yukawa couplings. To be consistent with a see-saw mechanism of neutrino masses, LH models should incorporate another mechanism to suppress the SM quadratic corrections at the ultraviolet cutoff $`\mathrm{\Lambda }10`$ TeV. In that case the role of LH ideas would be just, for example, to increase the scale of supersymmetry breaking in one order of magnitude. On the other hand, the possibility of very suppressed Yukawa couplings would imply that the LH model is embedded in a theory with extra dimensions or (its CFT dual) strongly coupled dynamics.
We have found a LH alternative that can explain the small size of neutrino masses with no need for a large scale nor extra dynamics. The lepton sector includes a gauge singlet $`n^c`$ per generation and has unsuppressed Yukawas, but it is free of one-loop quadratic corrections because of the same global symmetry as in the top-quark sector. This implies another $`SU(2)`$ singlet $`N`$ per generation, with $`n^c`$ and $`N^{}Nv/f\nu `$ combining into a massive field at the scale $`f`$ of global symmetry breaking. As long as lepton number is not broken, the SM neutrinos are massless. If $`L`$ is broken at a small scale $`M10^6`$ GeV (a mirror scale of the one in the see-saw mechanism) then the diagram in Fig. 3 introduces a term $`\varphi _2^{}\varphi _2^{}\mathrm{\Psi }_L\mathrm{\Psi }_L`$ and the SM neutrinos get the observed masses.
In just the SM with an extra singlet per generation this mechanism would not work: the small value of the singlet mass $`M`$ would not prevent the SM neutrinos to combine with the singlets and define Dirac fields of mass $`\lambda _\nu v`$ similar to the mass of quarks and charged leptons. However, the mechanism is naturally implemented in LH models, where the global symmetries imply new fields and allow at the loop level the necessary couplings.
We would like to thank Verónica Sanz for useful discussions. This work has been supported by MCYT (FPA2003-09298-C02-01) and Junta de Andalucía (FQM-101). J.L.P. acknowledges a FPU grant from MEC. |
warning/0506/hep-ph0506042.html | ar5iv | text | # 1 Introduction
## 1 Introduction
The mechanism of electroweak symmetry breaking (EWSB) remains the most pressing puzzle in elementary particle physics. Experimentally, this question will be addressed at the Large Hadron Collider (LHC). Theoretically, several interesting possibilities have been proposed. In this article, we will concentrate on the “Little Higgs” proposal . In this approach, the Higgs emerges as a pseudo-Nambu-Goldstone boson, whose properties are constrained by global symmetries. These global symmetries are not exact, and their breaking allows the Higgs to participate in non-derivative (i.e. gauge and Yukawa) interactions. At the same time, there is enough global symmetry left to ensure that the Higgs mass term vanishes at tree level, and is only logarithmically sensitive to the unknown short distance (ultraviolet, or UV) physics at the one-loop level. The usual quadratic sensitivity of the Higgs mass parameter on the UV physics first appears at two loops, and the incalculable UV effects remain subleading as long as the cutoff of the theory (the scale at which it becomes strongly coupled) is at or below about 10 TeV. With this requirement, the Higgs mass terms are dominated by the one-loop contribution from the top loops, which has the appropriate sign to trigger the electroweak symmetry breaking, and produces the Higgs vacuum expectation value (vev) of the right order of magnitude. Thus, Little Higgs theories provide an attractive explanation of EWSB.
The originally proposed implementations of the Little Higgs approach suffered from severe constraints from precision electroweak measurements , which could only be satisfied by finely tuning the model parameters. The most serious constraints resulted from the tree-level corrections to precision electroweak observables due to the exchanges of additional heavy gauge bosons present in the theories, as well as from the small but non-vanishing vev of an additional weak-triplet scalar field. Motivated by these constraints, several new implementations of the Little Higgs were proposed . Particularly interesting is the approach of Refs. , which introduces a discrete symmetry, dubbed “T parity” in analogy to R parity in the minimal supersymmetric standard model (MSSM). T parity explicitly forbids any tree-level contribution from the heavy gauge bosons to the observables involving only standard model (SM) particles as external states. It also forbids the interactions that induced the triplet vev. As a result, in T parity symmetric Little Higgs models, corrections to precision electroweak observables are generated exclusively at loop level. This implies that the constraints are generically weaker than in the tree-level case, and fine tuning can be avoided .
The main goal of this paper is to investigate the electroweak precision constraints on the models with T parity in more detail. We will concentrate on the T parity symmetric version of the Littlest Higgs (LH) model, based on an $`SU(5)/SO(5)`$ global symmetry breaking pattern . Some phenomenological aspects of this model have been analyzed in Ref. . The model possesses an attractive dark matter candidate, the T-odd partner of the hypercharge gauge boson, which has the correct relic abundance in certain regions of the parameter space. It also leads to an interesting set of signatures at the LHC; in particular, an excess of events with large missing transverse energy is expected. In this paper, we will compute the corrections to the properties of $`W/Z`$ bosons induced by the new particles present in the LH model, and perform a global fit to precision electroweak observables. We will show that a consistent fit can be obtained in a large region of the model parameter space, so that no significant fine tuning is required. We will also demonstrate that the LH model allows for consistent fits with values of the Higgs mass as large as 800 GeV, far in excess of the upper bound obtained within the standard model. Finally, we will show that there exists a non-vanishing overlap between the region allowed by precision electroweak fits and the region where the LH model provides all of the observed dark matter.
The rest of the paper is organized as follows. After briefly reviewing the $`SU(5)/SO(5)`$ Littlest Higgs model with T parity in Section 2, we will present the calculation of the corrections to precision electroweak observables in Section 3. In Section 4, we present the constraints on the parameter space of the model resulting from a global fit to precision electroweak observables. Section 5 contains our conclusions. A discussion of some aspects of the LH model in the renormalizable $`R_\xi `$ gauge, which we find useful in our calculations, is presented in the Appendix A.
## 2 The Model
In this section, we will review the LH model with T parity , emphasizing the features that will be important for the analysis of this paper.
### 2.1 Gauge-Scalar Sector
The Littlest Higgs model embeds the electroweak sector of the standard model in an $`SU(5)/SO(5)`$ non-linear sigma model (nl$`\sigma `$m). A global $`SU(5)`$ symmetry is broken to $`SO(5)`$ by the vev of an $`SU(5)`$ symmetric tensor $`\mathrm{\Sigma }`$ of the form
$$\mathrm{\Sigma }_0=\left(\begin{array}{ccccc}0& 0& 0& 1& 0\\ 0& 0& 0& 0& 1\\ 0& 0& 1& 0& 0\\ 1& 0& 0& 0& 0\\ 0& 1& 0& 0& 0\end{array}\right).$$
(2.1)
The low energy dynamics of the nl$`\sigma `$m is described in terms of the field
$$\mathrm{\Sigma }=e^{2i\mathrm{\Pi }/f}\mathrm{\Sigma }_0,$$
(2.2)
where $`\mathrm{\Pi }`$ is the “pion matrix” containing the Goldstone degrees of freedom, and $`f1`$ TeV is the nl$`\sigma `$m symmetry breaking scale, or “pion decay constant”. An $`[SU(2)\times U(1)]^2`$ subgroup of the global $`SU(5)`$ symmetry is gauged. The gauged generators have the form
$`Q_1^a=\left(\begin{array}{ccc}\sigma ^a/2& 0& 0\\ 0& 0& 0\\ 0& 0& 0\end{array}\right),`$ $`Y_1=\mathrm{diag}(3,3,2,2,2)/10,`$ (2.6)
$`Q_2^a=\left(\begin{array}{ccc}0& 0& 0\\ 0& 0& 0\\ 0& 0& \sigma ^a/2\end{array}\right),`$ $`Y_2=\mathrm{diag}(2,2,2,3,3)/10.`$ (2.10)
The kinetic term for the $`\mathrm{\Sigma }`$ field can be written as
$$_{\mathrm{kin}}=\frac{f^2}{8}\mathrm{Tr}D_\mu \mathrm{\Sigma }(D^\mu \mathrm{\Sigma })^{},$$
(2.11)
where
$$D_\mu \mathrm{\Sigma }=_\mu \mathrm{\Sigma }i\underset{j}{}\left[g_jW_j^a(Q_j^a\mathrm{\Sigma }+\mathrm{\Sigma }Q_j^{aT})+g_j^{}B_j(Y_j\mathrm{\Sigma }+\mathrm{\Sigma }Y_j)\right],$$
(2.12)
with $`j=1,2`$. Here, $`B_j`$ and $`W_j^a`$ are the $`U(1)_j`$ and $`SU(2)_j`$ gauge fields, respectively, and $`g_j^{}`$ and $`g_j`$ are the corresponding coupling constants. The vev $`\mathrm{\Sigma }_0`$ breaks the extended gauge group $`\left[SU(2)\times U(1)\right]^2`$ down to the diagonal subgroup, which is identified with the standard model electroweak group $`SU(2)_L\times U(1)_Y`$. The fourteen pions of the $`SU(5)/SO(5)`$ breaking decompose into representations of the electroweak gauge group as follows:
$$\mathrm{𝟏}_\mathrm{𝟎}\mathrm{𝟑}_\mathrm{𝟎}\mathrm{𝟐}_{\mathrm{𝟏}/\mathrm{𝟐}}\mathrm{𝟑}_{\pm \mathrm{𝟏}}.$$
(2.13)
We will denote the fields in the above four representations as $`\eta `$, $`\omega `$, $`H`$ and $`\varphi `$, respectively. The field $`H`$ has the appropriate quantum numbers to be identified with the SM Higgs; after EWSB, it can be decomposed as $`H=(i\pi ^+,\frac{v+h+i\pi ^0}{\sqrt{2}})^T`$, where $`v=246`$ GeV is the EWSB scale and $`h`$ is the physical Higgs field. Explicitly, the pion matrix in terms of these fields has the form
$$\mathrm{\Pi }=\left(\begin{array}{ccccc}\omega ^0/2\eta /\sqrt{20}& \omega ^+/\sqrt{2}& i\pi ^+/\sqrt{2}& i\varphi ^{++}& i\frac{\varphi ^+}{\sqrt{2}}\\ \omega ^{}/\sqrt{2}& \omega ^0/2\eta /\sqrt{20}& \frac{v+h+i\pi ^0}{2}& i\frac{\varphi ^+}{\sqrt{2}}& \frac{i\varphi ^0+\varphi _P^0}{\sqrt{2}}\\ i\pi ^{}/\sqrt{2}& (v+hi\pi ^0)/2& \sqrt{4/5}\eta & i\pi ^+/\sqrt{2}& (v+h+i\pi ^0)/2\\ i\varphi ^{}& i\frac{\varphi ^{}}{\sqrt{2}}& i\pi ^{}/\sqrt{2}& \omega ^0/2\eta /\sqrt{20}& \omega ^{}/\sqrt{2}\\ i\frac{\varphi ^{}}{\sqrt{2}}& \frac{i\varphi ^0+\varphi _P^0}{\sqrt{2}}& \frac{v+hi\pi ^0}{2}& \omega ^+/\sqrt{2}& \omega ^0/2\eta /\sqrt{20}\end{array}\right),$$
(2.14)
where the superscripts indicate the electric charge. The fields $`\eta `$ and $`\omega `$ are eaten<sup>1</sup><sup>1</sup>1In the LH model with T parity, the fields $`\eta `$ and $`\omega `$ mix with the field $`\varphi `$ at order $`(v/f)^2`$, and it is a linear combination of these that is eaten. See Appendix A for details. when the extended gauge group is broken down to $`SU(2)_L\times U(1)_Y`$, whereas the $`\pi `$ fields are absorbed by the standard model $`W/Z`$ bosons after EWSB. The fields $`h`$ and $`\varphi `$ remain in the spectrum. Including the EWSB effects, the vev of the $`\mathrm{\Sigma }`$ field has the form
$$\mathrm{\Sigma }=\left(\begin{array}{ccccc}0& 0& 0& 1& 0\\ 0& \frac{1}{2}(1c_v)& \frac{i}{\sqrt{2}}s_v& 0& \frac{1}{2}(1+c_v)\\ 0& \frac{i}{\sqrt{2}}s_v& c_v& 0& \frac{i}{\sqrt{2}}s_v\\ 1& 0& 0& 0& 0\\ 0& \frac{1}{2}(1+c_v)& \frac{i}{\sqrt{2}}s_v& 0& \frac{1}{2}(1c_v)\end{array}\right),$$
(2.15)
where
$$s_v=\mathrm{sin}\frac{\sqrt{2}v}{f},c_v=\mathrm{cos}\frac{\sqrt{2}v}{f}.$$
(2.16)
These formulas will prove very useful for analyzing the spectrum of the model.
The gauge generators are embedded in the $`SU(5)`$ is such a way that any given generator commutes with an $`SU(3)`$ subgroup of the $`SU(5)`$. This implies that if one pair of gauge couplings ($`g_1,g_1^{}`$ or $`g_2,g_2^{}`$) is set to zero, the Higgs field $`H`$ would be an exact Goldstone boson and, therefore, exactly massless. Thus, any diagram renormalizing the Higgs mass vanishes unless it involves at least two of the gauge couplings. At one loop, all diagrams satisfying this property are only logarithmically divergent: the “collective” symmetry breaking mechanism protects the Higgs mass from quadratic divergences. The first quadratic divergence appears at two loop level.
The original Littlest Higgs model described above turned out to be significantly constrained by precision electroweak observables . T parity, a discrete $`Z_2`$ symmetry, was introduced by Cheng and Low to avoid this difficulty, and it also provides a potential weak scale dark matter candidate. In the gauge sector, T parity is an automorphism of the gauge groups which exchanges the $`\left[SU(2)\times U(1)\right]_1`$ and $`\left[SU(2)\times U(1)\right]_2`$ gauge fields . The Lagrangian in Eq. (2.11) is invariant under this transformation provided that $`g_1=g_2`$ and $`g_1^{}=g_2^{}`$. In this case, the gauge boson mass eigenstates (before EWSB) have the simple form, $`W_\pm =(W_1\pm W_2)/\sqrt{2}`$, $`B_\pm =(B_1\pm B_2)/\sqrt{2}`$, where $`W_+`$ and $`B_+`$ are the standard model gauge bosons and are T-even, whereas $`W_{}`$ and $`B_{}`$ are the additional, heavy, T-odd states. (Typically, $`B_{}`$ is the lightest T-odd state, and plays the role of dark matter .) After EWSB, the T-even neutral states $`W_+^3`$ and $`B_+`$ mix to produce the SM $`Z`$ and the photon. Since they do not mix with the heavy T-odd states, the Weinberg angle is given by the SM relation, $`\mathrm{tan}\theta _w=g^{}/g`$, where $`g=g_{1,2}/\sqrt{2}`$ and $`g^{}=g_{1,2}^{}/\sqrt{2}`$ are the SM gauge couplings, and $`\rho =1`$ at tree level. As will be shown below, all the SM fermions are also T-even, implying that the $`W_{}`$ and $`B_{}`$ states generate no corrections to precision electroweak observables at tree level.
The transformation properties of the gauge fields under T parity and the structure of the Lagrangian (2.11) imply that T parity acts on the pion matrix as follows:
$$T:\mathrm{\Pi }\mathrm{\Omega }\mathrm{\Pi }\mathrm{\Omega },$$
(2.17)
where $`\mathrm{\Omega }=\mathrm{diag}(1,1,1,1,1)`$. This transformation law ensures that the complex $`SU(2)_L`$ triplet $`\varphi `$ is odd under T parity, while the Higgs doublet $`H`$ is even. The trilinear coupling of the form $`H^{}\varphi H`$ is therefore forbidden, and no triplet vev is generated. Eliminating this source of tree-level custodial $`SU(2)`$ violation further relaxes the precision electroweak constraints on the model.
### 2.2 Light Fermion Sector
In the original LH model, the fermion sector of the standard model remained unchanged with the exception of the third generation of quarks, where the top Yukawa coupling had to be modified to avoid the large quadratically divergent contribution to the Higgs mass from top loops. In the model with T parity, however, the standard model fermion doublet spectrum needs to be doubled to avoid compositeness constraints . For each lepton/quark doublet, two fermion doublets $`\psi _1(\mathrm{𝟐},\mathrm{𝟏})`$ and $`\psi _2(\mathrm{𝟏},\mathrm{𝟐})`$ are introduced. (The quantum numbers refer to representations under the $`SU(2)_1\times SU(2)_2`$ gauge symmetry.) These can be embedded in incomplete representations $`\mathrm{\Psi }_1,\mathrm{\Psi }_2`$ of the global $`SU(5)`$ symmetry. An additional set of fermions forming an $`SO(5)`$ multiplet $`\mathrm{\Psi }^c`$, transforming nonlinearly under the full $`SU(5)`$, is introduced to give mass to the extra fermions; the field content can be expressed as follows:
$$\begin{array}{ccc}\mathrm{\Psi }_1=\left(\begin{array}{c}\psi _1\\ 0\\ 0\end{array}\right),& \mathrm{\Psi }_2=\left(\begin{array}{c}0\\ 0\\ \psi _2\end{array}\right),& \mathrm{\Psi }^c=\left(\begin{array}{c}\psi ^c\\ \chi ^c\\ \stackrel{~}{\psi }^c\end{array}\right).\end{array}$$
(2.18)
These fields transform under the $`SU(5)`$ as follows:
$$\mathrm{\Psi }_1V^{}\mathrm{\Psi }_1,\mathrm{\Psi }_2V\mathrm{\Psi }_2,\mathrm{\Psi }^cU\mathrm{\Psi }^c,$$
(2.19)
where $`U`$ is the nonlinear transformation matrix defined in Refs. . The action of T parity on the multiplets takes
$$\mathrm{\Psi }_1\mathrm{\Sigma }_0\mathrm{\Psi }_2,\mathrm{\Psi }^c\mathrm{\Psi }^c.$$
(2.20)
These assignments allow a term in the Lagrangian of the form
$$\kappa f(\overline{\mathrm{\Psi }}_2\xi \mathrm{\Psi }^c+\overline{\mathrm{\Psi }}_1\mathrm{\Sigma }_0\mathrm{\Omega }\xi ^{}\mathrm{\Omega }\mathrm{\Psi }^c),$$
(2.21)
where $`\xi =\mathrm{exp}(i\mathrm{\Pi }/f)`$. This term gives a Dirac mass $`M_{}=\sqrt{2}\kappa f`$ to the T-odd linear combination of $`\psi _1`$ and $`\psi _2`$, $`\psi _{}=(\psi _1+\psi _2)/\sqrt{2}`$, together with $`\stackrel{~}{\psi }^c`$; the T-even linear combination, $`\psi _+=(\psi _1\psi _2)/\sqrt{2}`$, remains massless and is identified with the standard model lepton or quark doublet. To give Dirac masses to the remaining T-odd states $`\chi ^c`$ and $`\psi ^c`$, additional fermions with opposite gauge quantum numbers can be introduced .
To complete the discussion of the fermion sector, we introduce the usual SM set of the $`SU(2)_L`$-singlet leptons and quarks, which are T-even and can participate in the SM Yukawa interactions with $`\psi _+`$. The Yukawa interactions induce a one-loop quadratic divergence in the Higgs mass; however, the effect is numerically small except for the third generation of quarks. The Yukawa couplings of the third generation must be modified to incorporate the collective symmetry breaking pattern; this is discussed in the next subsection.
### 2.3 Top Sector
In order to avoid large one-loop quadratic divergences from the top sector, the $`\mathrm{\Psi }_1`$ and $`\mathrm{\Psi }_2`$ multiplets for the third generation must be completed to representations of the $`SU(3)_1`$ (“upper-left corner”) and $`SU(3)_2`$ (“lower-right corner”) subgroups of $`SU(5)`$. These multiplets are
$$\begin{array}{ccc}𝒬_1=\left(\begin{array}{c}q_1\\ U_{L1}\\ 0\end{array}\right),& 𝒬_2=\left(\begin{array}{c}0\\ U_{L2}\\ q_2\end{array}\right);& \end{array}$$
(2.22)
they obey the same transformation laws under T parity and the $`SU(5)`$ symmetry as do $`\mathrm{\Psi }_1`$ and $`\mathrm{\Psi }_2`$, see Eqs. (2.19) and (2.20). The quark doublets are embedded such that
$$q_i=\sigma _2\left(\begin{array}{c}u_{Li}\\ b_{Li}\end{array}\right).$$
(2.23)
In addition to the SM right-handed top quark field $`u_R`$, which is assumed to be T-even, the model contains two $`SU(2)_L`$-singlet fermions $`U_{R1}`$ and $`U_{R2}`$ of hypercharge 2/3, which transform under T parity as
$$U_{R1}U_{R2}.$$
(2.24)
The top Yukawa couplings arise from the Lagrangian of the form
$`_t`$ $`=`$ $`{\displaystyle \frac{1}{2\sqrt{2}}}\lambda _1fϵ_{ijk}ϵ_{xy}\left[(\overline{𝒬}_1)_i\mathrm{\Sigma }_{jx}\mathrm{\Sigma }_{ky}(\overline{𝒬}_2\mathrm{\Sigma }_0)_i\stackrel{~}{\mathrm{\Sigma }}_{jx}\stackrel{~}{\mathrm{\Sigma }}_{ky}\right]u_R`$ (2.25)
$`+\lambda _2f(\overline{U}_{L1}U_{R1}+\overline{U}_{L2}U_{R2})+\mathrm{h}.\mathrm{c}.`$
where $`\stackrel{~}{\mathrm{\Sigma }}=\mathrm{\Sigma }_0\mathrm{\Omega }\mathrm{\Sigma }^{}\mathrm{\Omega }\mathrm{\Sigma }_0`$ is the image of the $`\mathrm{\Sigma }`$ field under T parity, see Eq. (2.17), and the indices $`i,j,k`$ run from 1 to 3 whereas $`x,y=4,5`$. The T parity eigenstates are given by
$$q_\pm =\frac{1}{\sqrt{2}}(q_1q_2),U_{L\pm }=\frac{1}{\sqrt{2}}(U_{L1}U_{L2}),U_{R\pm }=\frac{1}{\sqrt{2}}(U_{R1}U_{R2}).$$
(2.26)
In terms of these eigenstates, Eq. (2.25) has the form
$$_\mathrm{m}^T=\lambda _1f\left[\frac{1}{2}(1+c_v)\overline{U}_{L+}+\frac{s_v}{\sqrt{2}}\overline{u}_{L+}\right]u_R+\lambda _2f\left(\overline{U}_{L+}U_{R+}+\overline{U}_LU_R\right)+\mathrm{h}.\mathrm{c}.$$
(2.27)
where we have used Eq. (2.15). The T-odd states $`U_L`$ and $`U_R`$ combine to form a Dirac fermion $`T_{}`$, with mass $`m_T_{}=\lambda _2f`$. The remaining T-odd states $`q_{}`$ receive a Dirac mass from the interaction in Eq. (2.21), and are assumed to be decoupled. The mass terms for the T-even states are diagonalized by defining
$`t_L`$ $`=`$ $`\mathrm{cos}\beta u_{L+}\mathrm{sin}\beta U_{L+},T_{L+}=\mathrm{sin}\beta u_{L+}+\mathrm{cos}\beta U_{L+},`$
$`t_R`$ $`=`$ $`\mathrm{cos}\alpha u_R\mathrm{sin}\alpha U_{R+},T_{R+}=\mathrm{sin}\alpha u_R+\mathrm{cos}\alpha U_{R+},`$ (2.28)
where $`t`$ is identified with the SM top and $`T_+`$ is its T-even heavy partner. The mixing angles are given by
$`\alpha `$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{tan}^1{\displaystyle \frac{4\lambda _1\lambda _2(1+c_v)}{4\lambda _2^2\lambda _1^2(2s_v^2+(1+c_v)^2)}},`$
$`\beta `$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{tan}^1{\displaystyle \frac{2\sqrt{2}\lambda _1^2s_v(1+c_v)}{4\lambda _2^2+(1+c_v)^2\lambda _1^22\lambda _1^2s_v}}.`$ (2.29)
To leading order in the $`v/f`$ expansion,
$$\mathrm{sin}\alpha =\frac{\lambda _1}{\sqrt{\lambda _1^2+\lambda _2^2}},\mathrm{sin}\beta =\frac{\lambda _1^2}{\lambda _1^2+\lambda _2^2}\frac{v}{f}.$$
(2.30)
The masses of the two T-even Dirac fermions are given by
$$m_{t,T_+}^2=f^2\mathrm{\Delta }\left(1\pm \sqrt{1\frac{\lambda _1^2\lambda _2^2s_v^2}{2\mathrm{\Delta }^2}}\right),$$
(2.31)
where
$$\mathrm{\Delta }=\frac{1}{2}\left(\lambda _2^2+\frac{\lambda _1^2}{2}(s_v^2+\frac{1}{2}(1+c_v)^2)\right).$$
(2.32)
To leading order in $`v/f`$,
$$m_t=\frac{\lambda _1\lambda _2v}{\sqrt{\lambda _1^2+\lambda _2^2}},m_{\mathrm{T}_+}=\sqrt{\lambda _1^2+\lambda _2^2}f.$$
(2.33)
It is interesting to note that the T-odd states do not participate in the cancellation of quadratic divergences in the top sector: the cancellation only involves loops of $`t`$ and $`T_+`$, and the details are identical to the LH model without T parity .
Using the above equations, it is straightforward to obtain the Feynman rules for the top sector of the LH model; we list the rules relevant for the calculations in this paper in Table 1.
## 3 Corrections to Precision Electroweak Observables
The introduction of T parity automatically eliminates the tree level electroweak precision constraints that plagued the original Littlest Higgs model: since the external states in all experimentally tested processes are T-even, no T-odd state can contribute to such processes at tree level. The only non-SM T-even state in our model, the heavy top $`T_+`$, can only contribute at tree level to observables involving the SM top quark<sup>2</sup><sup>2</sup>2Tree-level contributions of $`T_+`$ to other observables are suppressed by small off-diagonal CKM matrix elements., such as its couplings to $`W`$ and $`Z`$ bosons . At present, however, these couplings have not been measured experimentally, so no contraints arise at the tree level. At one-loop level, however, precision electroweak observables receive contributions from the $`T_+`$ as well as the T-odd particles. It is these contributions that determine the allowed parameter space of the Littlest Higgs model with T parity. We will evaluate the leading corrections in this section, and use them to perform a global fit to precision electroweak observables in Section 4.
In the SM, one-loop contributions to precision electroweak observables from the top sector, enhanced by powers of the top Yukawa coupling $`\lambda _t`$, dominate over contributions from the gauge and scalar sectors. We expect that the same hierarchy of effects will hold in the Littlest Higgs model, and our main focus will be on analyzing the effects of the top sector. However, we will also include the custodial-symmetry violating contributions from the gauge sector and the T-odd partners of light fermions, which become important in certain regions of the parameter space. In addition, we will show explicitly that the contributions from the complex scalar triplet, which were shown to be potentially important in the original LH model , completely decouple in the T parity symmetric case due to the absence of the triplet vev.
Before proceeding with the calculations, let us make the following comment. The nl$`\sigma `$m underlying the LH model is a non-renormalizable effective theory, valid up to a cutoff scale $`\mathrm{\Lambda }4\pi f`$. Every operator consistent with the symmetries of the low-energy theory will be generated at the cutoff scale $`\mathrm{\Lambda }`$, and will contribute to the precision electroweak observables. However, we do not include such operators in the fit. (The only exception we make is to include the leading operator contributing to the $`T`$ parameter, Eq. (3.40), since this parameter plays the most important role in constraining the model.) This is justified by the following considerations. First, while the contribution of the TeV-scale states that we will compute and the operator contributions that we will ignore are naively of the same order, $`v^2/(16\pi ^2f^2)v^2/\mathrm{\Lambda }^2`$, the former are logarithmically enhanced by a factor of $`\mathrm{log}(f^2/v^2)\mathrm{log}(\mathrm{\Lambda }^2/f^2)\mathrm{log}(4\pi )^25`$. Second, while a cancellation between the corrections computed below and the operator contribution is in principle possible if the UV physics produces an operator with a large coefficient, any significant change in the fits due to such a cancellation would represent a fine-tuning between the effects generated at two different energy scales, $`f`$ and $`\mathrm{\Lambda }`$.
### 3.1 Oblique Corrections
The largest corrections to precision electroweak observables in the LH model are induced by the one-loop diagrams involving the T-even $`T_+`$ quark shown in Fig. 1. These oblique corrections can be described in terms of the Peskin-Takeuchi $`S`$, $`T`$, and $`U`$ parameters . The calculation of these parameters is straightforward if somewhat tedious; the result is
$`S`$ $`=`$ $`{\displaystyle \frac{s_\beta ^2}{2\pi }}\left[\left({\displaystyle \frac{1}{3}}c_\beta ^2\right)\mathrm{log}x_t+{\displaystyle \frac{(1+x_t)^2}{(1x_t)^2}}+{\displaystyle \frac{2x_t^2(3x_t)\mathrm{log}x_t}{(1x_t)^3}}{\displaystyle \frac{8}{3}}\right],`$
$`T`$ $`=`$ $`{\displaystyle \frac{3}{16\pi }}{\displaystyle \frac{s_\beta ^2}{s_w^2c_w^2}}{\displaystyle \frac{m_t^2}{m_Z^2}}\left[{\displaystyle \frac{s_\beta ^2}{x_t}}1c_\beta ^2{\displaystyle \frac{2c_\beta ^2}{1x_t}}\mathrm{log}x_t\right],`$
$`U`$ $`=`$ $`{\displaystyle \frac{s_\beta ^2}{2\pi }}\left[s_\beta ^2\mathrm{log}x_t+{\displaystyle \frac{(1+x_t)^2}{(1x_t)^2}}+{\displaystyle \frac{2x_t^2(3x_t)\mathrm{log}x_t}{(1x_t)^3}}{\displaystyle \frac{8}{3}}\right],`$ (3.34)
where $`x_t=m_t^2/m_{\mathrm{T}_+}^2`$, $`s_\beta `$ is the sine of the left-handed $`tT_+`$ mixing angle given in Eq. (2.29), and $`s_w`$ is the sine of the Weinberg angle. In the limit when $`x_t1`$, these formulas simplify considerably and we obtain
$`S`$ $`=`$ $`{\displaystyle \frac{1}{3\pi }}\left({\displaystyle \frac{\lambda _1}{\lambda _2}}\right)^2{\displaystyle \frac{m_t^2}{m_{\mathrm{T}_+}^2}}\left[{\displaystyle \frac{5}{2}}+\mathrm{log}{\displaystyle \frac{m_{\mathrm{T}_+}^2}{m_t^2}}\right],`$
$`T`$ $`=`$ $`{\displaystyle \frac{3}{8\pi }}{\displaystyle \frac{1}{s_w^2c_w^2}}\left({\displaystyle \frac{\lambda _1}{\lambda _2}}\right)^2{\displaystyle \frac{m_t^4}{m_{\mathrm{T}_+}^2m_Z^2}}\left[\mathrm{log}{\displaystyle \frac{m_{\mathrm{T}_+}^2}{m_t^2}}1+{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{\lambda _1}{\lambda _2}}\right)^2\right],`$
$`U`$ $`=`$ $`{\displaystyle \frac{5}{6\pi }}\left({\displaystyle \frac{\lambda _1}{\lambda _2}}\right)^2{\displaystyle \frac{m_t^2}{m_{\mathrm{T}_+}^2}}.`$ (3.35)
The leading-order result for the $`T`$ parameter is in agreement with the previous analyses of LH models with and without T parity .
In contrast to $`T_+`$, the T-odd top partner $`T_{}`$ does not contribute to $`S,T`$ or $`U`$ since it is an $`SU(2)_L`$ singlet which does not mix with the SM top. However, $`T_{}`$ loops do affect precision electroweak observables at the level of $`(m_Z/m_T_{})^2`$ corrections which are not captured by the formalism of Peskin and Takeuchi. The corrections to the two most precisely measured observables, $`s_W(1m_W^2/m_Z^2)^{1/2}`$ and $`s_{}`$ (the value of the weak mixing angle implied from $`Z`$ decay asymmetries), are given by
$$s_{}^2s_0^2=\frac{2\alpha }{45\pi }\frac{c_w^2s_w^2}{c_w^2s_w^2}\frac{m_Z^2}{m_T_{}^2},s_W^2s_{}^2=0,$$
(3.36)
where $`s_0`$ is the reference value of the weak mixing angle inferred from
$$s_0^2(1s_0^2)=\frac{\pi \alpha }{\sqrt{2}G_Fm_Z^2}.$$
(3.37)
These corrections are very small, and we do not include them in the global fit performed in Section 4.
It is clear from Eqs. (3.35) that the $`T`$ parameter induced by the $`T_+`$ loops is about 20 times larger than the $`S`$ and the $`U`$ for the same model parameters, and therefore the constraints on the model are largely driven by the $`T`$ parameter. This parameter also receives a contribution from the gauge sector of the model. While in general subdominant, this correction becomes important when the $`tT_+`$ mixing is suppressed (namely when the ratio $`\lambda _1/\lambda _2`$ is small), and we will include it in the fit. This contribution arises from the custodial $`SU(2)`$-violating tree level mass splitting of the T-odd heavy $`W_H^3`$ and $`W_H^\pm `$ gauge bosons. Neglecting effects of order $`g^2`$, the mass splitting is given by
$$\mathrm{\Delta }M^2M^2(W_H^3)M^2(W_H^\pm )=\frac{g^2f^2}{4}\left(1c_v\right)^2\frac{1}{8}g^2\frac{v^4}{f^2}.$$
(3.38)
At one loop, this effect induces a contribution to the $`T`$ parameter :<sup>3</sup><sup>3</sup>3We find that the calculation of this contribution simplifies considerably in the Landau gauge, $`\xi =0`$.
$$T_{W_H}=\frac{9}{16\pi c_w^2s_w^2M_Z^2}\mathrm{\Delta }M^2\mathrm{log}\frac{\mathrm{\Lambda }^2}{M^2(W_H)}.$$
(3.39)
Note that the result is divergent, and depends on the UV cutoff of the theory $`\mathrm{\Lambda }`$. This should not be surprising since the theory we’re dealing with is non-renormalizable: indeed, the mass splitting in Eq. (3.38) comes from a dimension-6 operator of the form $`W^\mu W_\mu (H^{}H)^2`$, which appears when the $`\mathrm{\Sigma }`$ fields in Eq. (2.11) are expanded to order $`\mathrm{\Pi }^4`$. The UV divergence signals the presence of a “counterterm” operator of the form
$$_c=\delta _c\frac{g^2}{16\pi ^2}f^2\underset{i,a}{}\mathrm{Tr}\left[(Q_i^aD_\mu \mathrm{\Sigma })(Q_i^aD^\mu \mathrm{\Sigma })^{}\right],$$
(3.40)
where $`\delta _c`$ is an order-one coefficient whose exact value depends on the details of the UV physics. (The normalization of Eq. (3.40) is fixed by naive dimensional analysis .) Including the counterterm, the full contribution of the gauge sector to the $`T`$ parameter has the form
$$T=\frac{1}{4\pi s_w^2}\frac{v^2}{f^2}\left(\delta _c+\frac{9}{8}\mathrm{log}\frac{\mathrm{\Lambda }^2}{M^2(W_H)}\right)=\frac{1}{4\pi s_w^2}\frac{v^2}{f^2}\left(\delta _c+\frac{9}{4}\mathrm{log}\frac{4\pi }{g}\right).$$
(3.41)
where we have assumed $`\mathrm{\Lambda }=4\pi f`$ and used $`M(W_H)=gf`$. As expected, this contribution is parametrically subdominant to the correction from the top sector, Eq. (3.35), in the limit $`m_tm_Z`$. Furthermore, in agreement with the discussion at the beginning of this section, the effect of the operator (3.40) induced at the cutoff scale is subdominant, by a factor of $`1/\mathrm{log}(\mathrm{\Lambda }^2/M^2(W_H))0.2`$, compared to the calculable contribution in Eq. (3.39).
In Section 4, we will also be interested in the effects of varying the Higgs mass. To leading order, the Higgs contribution to the oblique parameters is given by
$`S`$ $`=`$ $`{\displaystyle \frac{1}{12\pi }}\mathrm{log}{\displaystyle \frac{m_h^2}{m_{h,\mathrm{ref}}^2}},`$
$`T`$ $`=`$ $`{\displaystyle \frac{3}{16\pi c_w^2}}\mathrm{log}{\displaystyle \frac{m_h^2}{m_{h,\mathrm{ref}}^2}},`$
$`U`$ $`=`$ $`0,`$ (3.42)
where $`m_{h,\mathrm{ref}}m_h`$ is the “reference” value of the Higgs mass used to obtain the SM predictions for precision electroweak observables<sup>4</sup><sup>4</sup>4It should be kept in mind that in the LH model, the Higgs couplings to the $`W/Z`$ bosons will receive corrections of order $`v/f`$, which have been neglected in Eq. (3.42). This will not affect any of the conclusions of our analysis in Section 4.. Interestingly, the negative contribution to the $`T`$ parameter from a heavy Higgs can be partially cancelled by the positive contribution from the $`T_+`$. As we will show below, this cancellation allows for a consistent fit to precision electroweak observables with the Higgs mass well above the upper bound obtained in the SM .
Finally, the LH model contains an additional T-odd $`SU(2)_L`$-triplet scalar field $`\varphi `$, with the mass
$$m_\varphi ^2\frac{2m_h^2f^2}{v^2}(1\mathrm{TeV})^2.$$
(3.43)
After EWSB, a mass splitting of order $`v^2/f`$ between various components of the triplet is generated, for example, by operators of the form $`H^{}\varphi \varphi ^{}H`$. Neglecting this mass splitting, the triplet contributions to $`S`$, $`T`$ and $`U`$ parameters vanish; keeping the terms of order $`(m_Z/m_\varphi )^2`$, its contributions to $`s_{}`$ and $`s_W`$ are given by
$$s_{}^2s_0^2=\frac{\alpha }{24\pi }\frac{s_w^2c_w^2}{c_w^2s_w^2}\frac{m_Z^2}{m_\varphi ^2},s_W^2s_{}^2=\frac{\alpha }{60\pi }\frac{m_W^2}{m_\varphi ^2}.$$
(3.44)
If the mass splitting is taken into account, non-zero contributions to the Peskin-Takeuchi parameters are induced; however, these effects are of order $`\mathrm{\Delta }m_\varphi ^2/m_\varphi ^2v^4/m_\varphi ^4`$, and are thus subleading to the corrections given in Eq. (3.44). We conclude that the effects of the triplet $`\varphi `$ on the precision electroweak observables in the LH model with T parity decouple with growing $`m_\varphi `$, and are negligible for $`m_\varphi `$ in its natural range, around 1 TeV. We will not include these effects in the global fit of Section 4.
### 3.2 Effects of the T-Odd Partners of Light Fermions
To implement T parity in the fermion sector of the LH model, it is necessary to introduce a T-odd fermion partner for each lepton/quark doublet of the SM (see Section 2.2). These particles are vector-like, and their effects on precision electroweak observables must decouple in the limit when their mass is taken to infinity. However, box diagrams involving the exchanges of Goldstone bosons $`\omega `$ and $`\eta `$, see Fig. 2, generate four-fermion operators whose coefficients increase if the mass of the T-odd fermions is increased while $`f`$ is kept fixed<sup>5</sup><sup>5</sup>5We are grateful to Thomas Gregoire for bringing this point to our attention.. This non-decoupling is easy to understand qualitatively: to increase the mass of the T-odd fermions, it is necessary to increase the Yukawa coupling $`\kappa `$ in Eq. (2.21), which in turn makes the four interaction vertices in the box diagrams stronger. Assuming that the couplings $`\kappa `$ are flavor-diagonal and flavor-independent, the generated operators have the form
$$𝒪_{4\mathrm{f}}=\frac{\kappa ^2}{128\pi ^2f^2}\overline{\psi }_L\gamma ^\mu \psi _L\overline{\psi }_L^{}\gamma _\mu \psi _L^{},$$
(3.45)
where $`\psi `$ and $`\psi ^{}`$ are (distinct) SM fermions, and we ignore the corrections of order $`g/\kappa `$. The experimental bounds on four-fermi interactions involving SM fields provide an upper bound on the T-odd fermion masses; the strongest constraint comes from the $`eedd`$ operator, whose coefficient is required to be smaller than $`2\pi /(26.4\mathrm{TeV})^2`$ . This yields
$$M_{\mathrm{TeV}}<4.8f_{\mathrm{TeV}}^2,$$
(3.46)
where $`M_{\mathrm{TeV}}`$ and $`f_{\mathrm{TeV}}`$ are the values of the T-odd fermion masses and the symmetry breaking scale, respectively, in TeV. In addition, a lower bound on the masses of the T-odd fermions can be obtained from non-observation of these particles at the Tevatron; in analogy with squarks of the MSSM, we expect the bound to be in the neighborhood of 250–300 GeV.
The contribution of each T-odd doublet to the $`T`$ parameter is given by
$$T_{\mathrm{T}\mathrm{odd}}=\frac{\kappa ^2}{192\pi ^2\alpha }\left(\frac{v}{f}\right)^2,$$
(3.47)
where we omit terms of order $`(v/f)^4`$ and higher. Note that, for a fixed value of $`f`$, this contribution increases with increasing T-odd fermion mass; Eq. (3.46) implies that
$$|T_{\mathrm{T}\mathrm{odd}}|<0.05,$$
(3.48)
independent of $`f`$. (Note that this bound relies on the assumption that the $`\kappa `$ couplings are flavor-independent.) Nevertheless, the T-odd fermions can have a noticeable effect on the precision electroweak fits due to a large number (twelve) of doublets in the SM; this will be illustrated in the next section.
### 3.3 $`Zb\overline{b}`$ Vertex Renormalization
In the SM, the largest non-oblique correction is the renormalization of the $`Zb\overline{b}`$ vertex by top quark loops. This effect is non-decoupling in the sense that it is proportional to the square of the top mass. This non-decoupling is most easily seen if the calculation is performed in the ’t Hooft-Feynman gauge . In this gauge, the non-decoupling part of the vertex correction comes purely from the diagrams involving the exchange of a Goldstone boson $`\pi ^\pm `$, since its couplings to the top and bottom quarks are enhanced by the top Yukawa $`\lambda _t`$. These diagrams are shown in Fig. 3. The diagrams involving the exchange of the gauge bosons are subdominant in the large-$`m_t`$ limit, and neglecting their contribution only induces an error of order $`(m_Z/m_t)^225`$% in the vertex correction calculation.
We have calculated the one-loop correction to the $`Zb\overline{b}`$ vertex in the LH model with T parity. We used the ’t Hooft-Feynman gauge (for a brief discussion of the $`R_\xi `$ gauges in the LH model, see Appendix A). As in the SM case, the diagrams involving $`\pi ^\pm `$ exchanges dominate in the large-$`m_t`$ limit, and we neglected all other contributions. (While a more precise calculation could be done, the effort would not be justified as the $`Zb\overline{b}`$ correction turns out to have only a small effect on precision electroweak fits.) These diagrams are of three kinds. First, the same diagrams as in the SM appear, but with the top coupling to the $`Z`$ modified according to Table 1. Second, all the diagrams in Fig. 3 also appear with the top replaced by the T-even heavy top partner $`T_+`$. Third, the “mixing” diagrams shown in Fig. 4 appear as a result of the mixing between $`t`$ and $`T_+`$. These diagrams can be easily calculated using the couplings given in Table 1 and in Eq. (A.61) of Appendix A. To leading order in the limit $`m_{\mathrm{T}_+}m_tm_W`$, the result is
$$\delta g_L^{b\overline{b}}=\frac{g}{c_w}\frac{\alpha }{8\pi s_w^2}\frac{m_t^4}{m_W^2m_{\mathrm{T}_+}^2}\frac{\lambda _1^2}{\lambda _2^2}\mathrm{log}\frac{m_{\mathrm{T}_+}^2}{m_t^2},$$
(3.49)
where $`\delta g_L^{b\overline{b}}`$ is the correction received by the $`Zb_L\overline{b}_L`$ vertex in the LH model in addition to the usual SM one-loop correction. It is interesting to note that this leading order contribution comes entirely from the mixing diagrams in Fig. 4. The correction to the $`Zb_R\overline{b}_R`$ vertex is negligible since it is not enhanced by the top Yukawa coupling. Note also that the correction in Eq. (3.49) does not have the correct sign to alleviate the well-known deviation of the measured value of the forward-backward asymmetry in $`Zb\overline{b}`$ decays from the SM prediction inferred from the other precision electroweak observables .
A calculation of $`Zb\overline{b}`$ in a general theory containing an extra heavy quark that mixes with $`t`$ has been carried out in Ref. . Accounting for the fact that $`U_+`$ is a vector isosinglet and including the appropriate mixing specific to the LH model, Eq. (3.49) agrees with the results of this analysis in the limit $`m_{\mathrm{T}_+}m_tm_W`$. The results of Ref. are more general, valid for arbitrary values of $`m_{\mathrm{T}_+}`$ and $`m_t`$. However, we find that using these expressions instead of Eq. (3.49) does not lead to noticeable changes in the global fits performed in Section 4.
With the assumption of flavor-diagonal and flavor-independent Yukawa couplings $`\kappa `$ made in Section 3.2, the one-loop vertex corrections due to loops of T-odd fermions are flavor-universal, and can therefore be absorbed in the redefinitions of gauge couplings. They will not induce an observable shift in $`Zb\overline{b}`$ couplings.
## 4 Constraints on the Littlest Higgs Parameter Space
To obtain constraints on the parameter space of the LH model with T parity, we have performed a global fit to precision electroweak observables, including the LH contributions evaluated in the previous section. The LH contributions are parametrized by two dimensionless numbers, $`R=\lambda _1/\lambda _2`$ and $`\delta _c`$, and the symmetry breaking scale $`f`$. In the fit, we have used the values of the 21 $`Z`$ pole and low-energy observables listed in Ref. ; the equations expressing the shifts in these observables in terms of the oblique parameters and $`\delta g_L^{b\overline{b}}`$ are given in Ref. . We take the top mass to be 176.9 GeV , and do not include the uncertainty associated with the top mass. In each constraint plot, we draw the $`95`$, $`99`$, and $`99.9\%`$ confidence level contours in the context of a $`\chi ^2`$ analysis with two degrees of freedom<sup>6</sup><sup>6</sup>6It is important to note that changing the assumed number of degrees of freedom can strongly affect the positions of the contours; this is equivalent to modifying the priors that enter into the fit . A complete Bayesian analysis taking into account a variety of different priors for the model parameters is beyond the scope of this paper, but it would be straightforward to perform such an analysis using the formulas provided in Section 3..
In the first part of the analysis, we have fixed the Higgs mass at its reference value, $`m_{h,\mathrm{ref}}=113`$ GeV. In Fig. 5, we plot the constraints in the $`fR`$ plane, assuming $`\delta _c=0`$. In Fig. 6, we fix $`R=1`$ and plot the constraints in the $`f\delta _c`$ plane, neglecting the T-odd fermion contribution. It is clear that a large part of the parameter space is consistent with precision electroweak constraints, including regions where the symmetry breaking scale $`f`$ is as low as 500 GeV. (In these regions, a partial cancellation between the top and gauge sector contributions to the $`T`$ parameter takes place.) In some cases, we have even obtained consistent fits for values of $`f`$ as low as 350 GeV. However, since our analysis neglects all higher-derivative operators generated at the scale $`\mathrm{\Lambda }4\pi f`$, which can contribute significantly to precision electroweak observables when $`\mathrm{\Lambda }<5`$ TeV, we estimate that the fits cannot be trusted for $`f<400`$ GeV, and do not show that part of the parameter space in the plots.
As has been shown in Section 3.1, top sector loops in the LH model provide a sizable, positive contribution to the $`T`$ parameter. This raises an interesting possibility: since the contribution of a heavy SM Higgs to the $`T`$ parameter is negative, it is possible that these two effects partially cancel<sup>7</sup><sup>7</sup>7A consistent fit with a heavy Higgs can also be obtained in the Littlest Higgs model without T parity, where a positive correction to the $`T`$ parameter is generated at tree level; however, this requires a rather high value of $`f`$, of order $`5`$ TeV . A similar cancellation of the heavy Higgs and new physics contributions to $`T`$ also occurs in top seesaw models ; see Ref. . We are grateful to Bogdan Dobrescu for bringing this paper to our attention., and a consistent fit is obtained for $`m_h`$ far in excess of the usual SM upper bound, currently about 250 GeV. This possibility is illustrated in Fig. 7, where we fix $`R=2`$, $`\delta _c=0`$, and plot the constraints in the $`fm_h`$ plane. Remarkably, values of $`m_h`$ as high as 800 GeV are allowed at 95% confidence level. (Note that the approximation made in Eq. (3.42), where the corrections of order $`v/f`$ in the Higgs contribution to the oblique parameters have been neglected, is justified in the region of interest, since $`f`$ is still of order 1 TeV.) Thus, the LH model provides an explicit, well-motivated example of a theory in which the SM upper bound on the Higgs mass is avoided. Moreover, from the point of view of fine tuning in the Higgs potential, the high values of $`m_h`$ are more natural in the context of this model . For example, let us use the ratio of the one-loop top contribution to $`m_h^2`$ to the full $`m_h^2`$,
$$F=\frac{3\lambda _t^2m_{\mathrm{T}_+}^2}{4\pi ^2m_h^2}\mathrm{log}\frac{\mathrm{\Lambda }^2}{m_{\mathrm{T}_+}^2},$$
(4.50)
as a quantitative measure of the fine tuning. (Larger values of $`F`$ correspond to higher degree of fine tuning.) Plotting the contours of constant $`F`$ indicates that, in the region of the parameter space consistent with precision electroweak constraints, the degree of fine tuning increases with the decreasing Higgs mass. Large values of $`m_h`$ are clearly preferred from the point of view of naturalness in the Higgs potential.
If T parity is an exact symmetry (including the theory completing the description above the scale $`\mathrm{\Lambda }`$), the lightest T-odd particle (LTP) is stable. Generically, the LTP is the T-odd partner of the hypercharge gauge boson, which is electrically neutral and can play the role of WIMP dark matter. The LTP relic density has been computed in Ref. , and a region in the parameter space where the LTP can account for all of the observed dark matter has been identified. In Fig. 8, the 2 sigma contours on the dark matter relic density are superposed over a plot of the precision electroweak constraints where $`R=2`$, $`\delta _c=0`$, and $`m_h`$ and $`f`$ are allowed to vary. There is a region of the allowed parameter space in which the LTP can account for all of the dark matter.
In Figs. 58, the contribution of the T-odd fermions to the $`T`$ parameter is neglected. This approximation is justified as long as the T-odd fermions are sufficiently light: for example, for T-odd fermion mass of 300 GeV, their total contribution to the $`T`$ parameter is very small, and does not have any noticeable effect on the fits. On the other hand, heavier T-odd fermions can have a substantial effect. This is illustrated in Fig. 9, where the T-odd fermion contribution has been assumed to have the maximal size consistent with the constraint from four-fermi interactions, Eq. (3.48). (This corresponds to the T-odd fermion masses saturating the upper bound in Eq. (3.46).) While the constraints in this case are more severe, consistent fits can still be obtained for values of $`f`$ below 1 TeV.
## 5 Conclusions
In this paper we have calculated the dominant corrections to the precision electroweak observables at the one-loop level in the Littlest Higgs model with T parity . We performed a global fit to the precision electroweak observables and found that a large part of the model parameter space is consistent with data. In particular, a consistent fit can be obtained for values of the nl$`\sigma `$m symmetry breaking scale $`f`$ as low as 500 GeV. Furthermore, we found that the LH model can fit the data for values of the Higgs mass far in excess of the SM upper bound, due to the possibility of a partial cancellation between the contributions to the $`T`$ parameter from Higgs loops and new physics. Combining our results with those of Ref. , we found that there are regions of parameter space allowed by precision electroweak constraints where the lightest T-odd particle can account for all of the observed dark matter.
We have argued that the corrections to low energy observables in the LH model are dominated by the top sector, and our analysis was primarily focused on those contributions. It would be interesting to perform a more detailed analysis of the effects from the gauge and scalar sectors; however, we do not expect these effects to substantially modify our conclusions. The analysis of the T-odd fermion sector in this paper relied on rather restrictive simplifying assumptions: in particular, the Yukawa couplings in the T-odd sector were assumed to be flavor-diagonal and flavor-independent. A possible non-trivial flavor structure of their couplings could have interesting experimental consequences. Moreover, these fermions should be sufficiently light to be pair-produced at the LHC, or even at the Tevatron. It is therefore important to analyze that sector of the model in more detail.
In conclusion we find that the Littlest Higgs model with T parity is only weakly constrained by precision electroweak data, and provides a viable alternative for physics at the TeV scale. Apart from being theoretically attractive, the model has several features that are of interest for planning future experiments. Two examples are the possibility of a relatively heavy Higgs, as discussed in this paper, and the similarity of many of the collider signatures of this model to the benchmark SUSY signatures, which will inevitably complicate the LHC analysis . We hope that our analysis, which explicitly demonstrates the viability of the LH model, will open the door for further detailed studies of its collider phenomenology.
## Acknowledgements
We would like to thank Bogdan Dobrescu and Ian Low for useful discussions related to this work. We would also like to acknowledge helpful correspondence with Thomas Gregoire. This work is supported by the National Science Foundation under grant PHY-0355005.
## Appendix A Renormalizable Gauges for the Littlest Higgs with T Parity
While the higher-order corrections for observable quantities in gauge theories must be gauge independent, an appropriate gauge choice can greatly reduce the complexity of a loop calculation, and make the underlying physics more transparent. This is especially important in the case of spontaneously broken gauge symmetries. While many issues in these theories are most easily understood in the unitary gauge, this gauge is ill-suited for loop calculations, leading to complicated intermediate expressions, and, in some cases, ambiguous answers<sup>8</sup><sup>8</sup>8A well-known example of such an ambiguity appears in the calculation of the $`W`$ boson contribution to the anomalous magnetic moment of the muon .. Experience with the SM radiative correction calculations indicates that it is best to use the renormalizable, or $`R_\xi `$, gauges; a special case of $`\xi =1`$, the ’t Hooft-Feynman gauge, is especially useful. In this paper, we have used this gauge to calculate the $`Zb\overline{b}`$ coupling shift, see Section 3.3. Since only unitary-gauge Feynman rules have appeared in the literature so far for the LH models , we will briefly discuss the $`R_\xi `$ gauges for the LH model with T parity in this Appendix. We will focus on the charged gauge boson sector; the analysis of the neutral sector is similar. Even though the calculations in the paper do not require it, in this Appendix we will keep all correction of order $`(v/f)^2`$, since an interesting effect of $`\omega \varphi `$ mixing first appears at that order.
The charged gauge boson mass matrix follows from Eq. (2.11); to order $`ϵ^2`$ (where $`ϵ=v/f`$) it has the form
$$M^2=f^2\left(\begin{array}{cc}g_1^2& g_1g_2(1ϵ^2/4)\\ g_1g_2(1ϵ^2/4)& g_2^2\end{array}\right).$$
(A.51)
Diagonalizing this matrix results in the mass eigenstates
$$W_L=c_0W_1+s_0W_2,W_H=s_0W_1+c_0W_2.$$
(A.52)
In the LH model with T parity, the gauge couplings are set equal, $`g_1=g_2=\sqrt{2}g`$ and the mixing angle is given by $`s_0=c_0=1/\sqrt{2}`$. The charged gauge boson mass eigenvalues are then
$$M_H^2=g^2f^2\left[1\frac{1}{4}ϵ^2+\mathrm{}\right],M_L^2=\frac{g^2v^2}{4}\left[1\frac{1}{6}ϵ^2+\mathrm{}\right].$$
(A.53)
Spontaneous breaking of the gauge symmetries leads to the kinetic mixing between the gauge bosons and the Goldstone boson fields in Eq. (2.14). In the mass eigenbasis, the mixing terms have the form
$`_{W\pi }`$ $`=`$ $`M_HW_H^\mu \left[_\mu \omega ^+\left(1{\displaystyle \frac{1}{24}}ϵ^2\right){\displaystyle \frac{i}{6}}ϵ^2_\mu \varphi ^+\right]`$ (A.54)
$`+M_LW_L^{}_\mu \pi ^+\left(1{\displaystyle \frac{1}{12}}ϵ^2\right)+\mathrm{h}.\mathrm{c}.`$
At order $`ϵ^2`$, the Goldstone boson fields in Eq. (2.14) are not canonically normalized. To canonically normalize the Goldstone fields, we perform the following rescaling:
$`\pi ^\pm `$ $``$ $`\pi ^\pm \left(1+{\displaystyle \frac{1}{12}}ϵ^2\right),`$
$`\omega ^\pm `$ $``$ $`\omega ^\pm \left(1+{\displaystyle \frac{1}{24}}ϵ^2\right),`$
$`\varphi ^\pm `$ $``$ $`\varphi ^\pm \left(1+{\displaystyle \frac{1}{24}}ϵ^2\right).`$ (A.55)
After this redefinition, there are still kinetic mixing terms involving the $`\omega `$ fields and the complex triplet, $`\varphi `$:
$$_{\mathrm{kin}}=_\mu \omega ^+^\mu \omega ^{}+_\mu \varphi ^+^\mu \varphi ^{}+\frac{i}{12}ϵ^2(_\mu \omega ^+^\mu \varphi ^{}_\mu \varphi ^+^\mu \omega ^{}).$$
(A.56)
These terms are diagonalized with the redefinition
$`\omega ^\pm =\omega ^\pm {\displaystyle \frac{i}{24}}ϵ^2\varphi ^\pm ,`$
$`\varphi ^\pm =\varphi ^\pm \pm {\displaystyle \frac{i}{24}}ϵ^2\omega ^\pm .`$ (A.57)
In terms of these new canonically normalized fields, the gauge boson-Goldstone mixing terms are given by<sup>9</sup><sup>9</sup>9Note that the normalizations of the fields in the definition of the Goldstone boson matrix, Eq. (2.14), have been chosen so that the mixing term has the simple form in Eq. (A.58).
$`_{W\pi }`$ $`=`$ $`M_HW_H^\mu \left[_\mu \omega ^+{\displaystyle \frac{i}{8}}ϵ^2_\mu \varphi ^+\right]`$ (A.58)
$`+M_LW_L^{}_\mu \pi ^++\mathrm{h}.\mathrm{c}.`$
A final rotation which leaves the kinetic terms diagonal,
$`\omega ^{\prime \prime \pm }=\omega ^\pm {\displaystyle \frac{i}{8}}ϵ^2\varphi ^\pm ,`$
$`\varphi ^{\prime \prime \pm }=\varphi ^\pm {\displaystyle \frac{i}{8}}ϵ^2\omega ^\pm ,`$ (A.59)
identifies $`\omega ^{\prime \prime \pm }`$ as the combination of Goldstones eaten by the heavy gauge bosons, and $`\varphi ^{\prime \prime \pm }`$ as the uneaten combination. A similar, but algebraically more involved, analysis carries through for the massive neutral gauge bosons.
Following the usual logic of renormalizable gauges, we add a gauge-fixing term which, after integration by parts, cancels the mixing terms:
$$\mathrm{\Delta }=\frac{1}{2\xi _L}\left|_\mu W_{L\mu }^\pm +M_L\xi _L\pi ^\pm \right|^2+\frac{1}{2\xi _H}\left|_\mu W_H^\pm +M_H\xi _H\omega ^{\prime \prime \pm }\right|^2.$$
(A.60)
The mass eigenstates in the eaten Goldstone sector are $`\pi ^\pm `$, with mass $`\sqrt{\xi _L}M_L`$, and $`\omega ^{\prime \prime \pm }`$, with mass $`\sqrt{\xi _H}M_H`$. (We have used the ’t Hooft–Feynman gauge, $`\xi _L=1`$, in Section 3.3.) Note that the $`\pi `$ fields do not mix with $`\omega `$ and $`\varphi `$ at any order in $`v/f`$, since such mixing is forbidden by T parity. The situation would be considerably more involved in the case of the Littlest Higgs without T parity.
Given the exact identification of the $`\pi ^\pm `$ fields in Eq. (2.14) with the lighter mass eigenstate in the $`R_\xi `$ gauges, it is straightforward to obtain the $`bt\pi `$ and $`bT\pi `$ vertices required in the calculation of Section 3.3. Expanding the $`\mathrm{\Sigma }`$ matrices in Eq. (2.25) to linear order in $`\mathrm{\Pi }`$, and using Eqs. (2.26), (2.28) to transform to the mass eigenbasis for the top sector, we obtain
$$i\sqrt{2}\lambda _1\overline{b}_{L+}\pi ^{}(\mathrm{cos}\alpha t_R+\mathrm{sin}\alpha T_{R+})+\mathrm{h}.\mathrm{c}.=i\lambda _t\overline{b}_{L+}\pi ^{}(t_R+\frac{\lambda _1}{\lambda _2}T_{R+})+\mathrm{h}.\mathrm{c}.$$
(A.61)
where $`\lambda _t`$ is the SM Yukawa coupling, and $`b_+`$ is identified with the SM $`b`$ quark. Note that the couplings involving the $`b_+`$ and any one of the T-odd Goldstone bosons, $`\varphi `$, $`\omega `$ or $`\eta `$, vanish due to the structure of the Lagrangian (2.25) and the fact that the field $`u_R`$ is T-even. |
warning/0506/math0506154.html | ar5iv | text | # SKEW DERIVATIONS AND DEFORMATIONS OF A FAMILY OF GROUP CROSSED PRODUCTS
## 1. Introduction
Deformations of a polynomial algebra, such as the Weyl algebra or functions on quantum affine space, may be expressed by formulas involving derivations of the polynomial algebra. These formulas are power series in an indeterminate with coefficients in the universal enveloping algebra of the Lie algebra of derivations. There are generalizations of such deformations to other types of algebras, such as functions on a manifold or orbifold, that are of current interest.
In this note we give a new generalization of the formulas themselves, and apply it to crossed products of polynomial algebras with groups of linear automorphisms. These group crossed products are of interest in geometry due to their relationship with corresponding orbifolds. Particular deformations of such crossed products, called graded Hecke algebras, were defined by Drinfel’d . These deformations have been studied by many authors, for example for crossed products with real reflection groups, see , with complex reflection groups, see , and with symplectic reflection groups, see . For these crossed product algebras, the universal enveloping algebra of the Lie algebra of derivations does not capture all the known deformations. Instead we derive a deformation formula from the action of a Hopf algebra under some hypotheses, recovering more of these known deformations as well as some new ones.
We use the theory developed by Giaquinto and Zhang of a universal deformation formula based on a bialgebra $`B`$ , extending earlier such formulas based on universal enveloping algebras of Lie algebras. Such a formula is universal in the sense that it applies to any $`B`$-module algebra to yield a formal deformation. Known examples include formulas based on universal enveloping algebras of Lie algebras (see examples and references in ) and a formula based on a small noncocommutative bialgebra . In Section 3 we generalize the formula in . Our universal deformation formula is based on a bialgebra generated by skew-primitive and group-like elements, and depends on a parameter $`q`$. The bialgebra and formula were discovered in the generic case by Giaquinto and Zhang, but were not published . We state their formula and modify it to include the case where $`q`$ is a root of unity (Theorem 3.3). The case $`q=1`$ is \[4, Lemma 6.2\], and as a first new example, we show that the smallest Taft algebra may be deformed by applying this earlier formula (Example 3.4). This is one of a series of algebras defined by quiver and relations whose deformations were given by Cibils , and we recover one of his deformations. We end Section 3 with more details on the structure of the bialgebras involved, briefly reviewing Kharchenko’s construction of a Hopf algebra of automorphisms and skew derivations of an algebra .
We apply our universal deformation formula in Section 4 to deform some crossed products of polynomial algebras with groups (Corollary 4.8), generalizing \[4, Example 6.3\] in which the group was $`/2\times /2`$. We use Hochschild cohomology to prove that the resulting deformations are nontrivial, by showing that their associated infinitesimals are not coboundaries. We apply the results of to show how our deformation formula leads, in special cases, to (twisted) graded Hecke algebras (Example 4.13). Due to restrictive hypotheses, our formula does not give rise to very many of the examples of (twisted) graded Hecke algebras in . However it does give an infinite series of universal deformation formulas based on noncocommutative Hopf algebras, as well as algebras thereby undergoing nontrivial deformation. It also proves existence of more general deformations of certain crossed products than those that are the (twisted) graded Hecke algebras.
In Section 5, we give two small examples for which our universal deformation formula nearly provides a universal deformation of the crossed product algebra in the other sense of the word universal. That is all possible nonclassical deformations of the algebra are parametrized by the formula, where by nonclassical we mean those not arising from deformations of the underlying polynomial algebra itself.
For completeness, we include an appendix in which the Hochschild cohomology of the relevant crossed product algebra is computed. In the case of a trivial twisting two-cocycle associated to the group, this was done elegantly by Farinati , and by Ginzburg and Kaledin in a more general geometric setting. Their results are easily generalized to crossed products with twisting cocycles, however in Section 4 we need some details from a more direct calculation. We give an algebraic computation similar to that in where the crossed product was taken with a Weyl algebra instead of a polynomial algebra (and the group is symplectic). We thank R.-O. Buchweitz for explaining this computation of Hochschild cohomology to us.
We also owe many thanks to A. Giaquinto and J. Zhang for sharing their unpublished universal deformation formula with us, and especially to A. Giaquinto from whom we learned algebraic deformation theory. We thank C. Cibils and J. Stasheff for many comments on earlier versions of this paper.
We will work over the complex numbers, although the definitions make sense more generally. Unless otherwise indicated, $`=_{}`$.
## 2. Definitions
Let $`S`$ be a $``$-algebra. Denote by $`\mathrm{Aut}_{}(S)`$ the group of all $``$-algebra automorphisms of $`S`$ that preserve the multiplicative identity. Let $`g,h\mathrm{Aut}_{}(S)`$. A $`g,h`$-skew derivation of $`S`$ is a $``$-linear function $`D:SS`$ such that
$$D(rs)=D(r)g(s)+h(r)D(s)$$
for all $`r,sS`$. If $`g=h=1`$ (the identity automorphism), then $`D`$ is simply a derivation of $`S`$.
We will be interested in skew derivations of a crossed product algebra which we define next. For more details on group crossed products, see . Let $`G`$ be any subgroup of $`\mathrm{Aut}_{}(S)`$. Let $`\alpha :G\times G^\times `$ be a two-cocycle, that is a function satisfying
(2.1)
$$\alpha (g,h)\alpha (gh,k)=\alpha (g,hk)\alpha (h,k)$$
for all $`g,h,kG`$. The crossed product ring $`S\mathrm{\#}_\alpha G`$ is $`SG`$ as a vector space, with multiplication
$$(rg)(sh)=\alpha (g,h)rg(s)gh$$
for all $`r,sS`$ and $`g,hG`$. This product is associative as $`\alpha `$ is a two-cocycle. We say $`\alpha `$ is a coboundary if there is some function $`\beta :G^\times `$ such that $`\alpha (g,h)=\beta (g)\beta (h)\beta (gh)^1`$ for all $`g,hG`$. The set of two-cocycles modulo coboundaries forms an abelian group under pointwise multiplication, that is $`(\alpha \alpha ^{})(g,h)=\alpha (g,h)\alpha ^{}(g,h)`$ for all $`g,hG`$. The crossed product algebras $`S\mathrm{\#}_\alpha G`$ and $`S\mathrm{\#}_\alpha ^{}G`$ are isomorphic if $`\alpha ^{}=\alpha \beta `$ for some coboundary $`\beta `$ (that is $`\alpha `$ and $`\alpha ^{}`$ are cohomologous).
We will abbreviate the element $`rg`$ of $`S\mathrm{\#}_\alpha G`$ by $`r\overline{g}`$. We will assume that $`\alpha `$ is normalized so that $`\alpha (1,g)=\alpha (g,1)=1`$ for all $`gG`$. Thus $`\overline{1}`$ is the multiplicative identity of $`S\mathrm{\#}_\alpha G`$, and it also follows from this and (2.1) that $`\alpha (g,g^1)=\alpha (g^1,g)`$ for all $`gG`$. The action of $`G`$ on $`S`$ extends to an inner action on $`S\mathrm{\#}_\alpha G`$, with $`g(a)=\overline{g}a(\overline{g})^1`$ for all $`gG`$, $`aS\mathrm{\#}_\alpha G`$, where $`(\overline{g})^1=\alpha ^1(g,g^1)\overline{g^1}=\alpha ^1(g^1,g)\overline{g^1}`$.
Now let $`t`$ be an indeterminate. A formal deformation of a $``$-algebra $`A`$ (for example $`A=S\mathrm{\#}_\alpha G`$) is an associative algebra $`A[[t]]=[[t]]A`$ over formal power series $`[[t]]`$ with multiplication
(2.2)
$$ab=ab+\mu _1(ab)t+\mu _2(ab)t^2+\mathrm{}$$
for all $`a,bA`$, where $`ab`$ denotes the product in $`A`$ and the $`\mu _i:AAA`$ are $``$-linear maps extended to be $`[[t]]`$-linear. Associativity of $`A[[t]]`$ implies that $`\mu _1`$ is a Hochschild two-cocycle, that is
(2.3)
$$\mu _1(ab)c+\mu _1(abc)=\mu _1(abc)+a\mu _1(bc)$$
for all $`a,b,cA`$, as well as further conditions on the $`\mu _i`$, $`i1`$. Thus a Hochschild two-cocycle $`\mu _1`$ is the first step towards a formal deformation, and it is called the infinitesimal of the deformation. In general it is difficult to determine whether a given $`\mu _1`$ lifts to a formal deformation of $`A`$. Hochschild cohomology is defined in the appendix; for more details on Hochschild cohomology and deformations of algebras, see .
One way in which to obtain a formal deformation of $`A`$ is through the action of a bialgebra on $`A`$. A bialgebra over $``$ is an associative $``$-algebra $`B`$ with $``$-algebra maps $`\mathrm{\Delta }:BBB`$ and $`\epsilon :B`$ such that $`(\mathrm{\Delta }\mathrm{id})\mathrm{\Delta }=\mathrm{\Delta }(\mathrm{id}\mathrm{\Delta })`$ and $`(\epsilon \mathrm{id})\mathrm{\Delta }=\mathrm{id}=(\mathrm{id}\epsilon )\mathrm{\Delta }`$. We will use the standard notation $`\mathrm{\Delta }(b)=b_1b_2`$ for all $`bB`$, where the subscripts are merely place-holders. A Hopf algebra is a bialgebra $`H`$ with a $``$-linear map $`S:HH`$ such that $`(Sh_1)h_2=\epsilon (h)=h_1(Sh_2)`$ for all $`hH`$. For details on bialgebras and Hopf algebras, see .
A universal deformation formula (or UDF) based on a bialgebra $`B`$ is an element $`F(BB)[[t]]`$ of the form $`F=11+tF_1+t^2F_2+\mathrm{}`$ with each $`F_iBB`$, satisfying
(2.4)
$$(\epsilon \mathrm{id})(F)=11=(\mathrm{id}\epsilon )(F)$$
$$\text{and }[(\mathrm{\Delta }\mathrm{id})(F)](F1)=[(\mathrm{id}\mathrm{\Delta })(F)](1F),$$
where $`\mathrm{id}`$ denotes the identity map. These relations (2.4) are similar to some of the defining relations for an $`R`$-matrix, and in fact both the inverse $`R^1`$ and the transpose $`R_{21}`$ of an $`R`$-matrix satisfy (2.4) (see ).
Suppose $`A`$ is a $``$-algebra which is also a left $`B`$-module. Then $`A`$ is a left $`B`$-module algebra if
(2.5)
$$h(ab)=(h_1a)(h_2b)\text{ and }h1=\epsilon (h)1$$
for all $`a,bA`$ and $`hB`$. This may be extended to a $`[[t]]`$-linear action of $`B[[t]]`$ by extending the scalars for $`A`$ to $`[[t]]`$. Let $`m:AAA`$ denote multiplication in $`A`$, extended to be $`[[t]]`$-linear. The following proposition combines Theorem 1.3 and Definition 1.13 of , and we sketch a proof for completeness.
###### Proposition 2.6 (Giaquinto-Zhang).
Let $`B`$ be a bialgebra, $`A`$ a left $`B`$-module algebra, and $`F`$ a universal deformation formula based on $`B`$. There is a formal deformation of $`A`$ given by $`ab=(mF)(ab)`$ for all $`a,bA`$.
###### Proof.
The format of $`F`$ as a power series in $`t`$ implies that $`ab=(mF)(ab)`$ takes the form (2.2). Associativity of $``$ follows from the second relation in (2.4) and the first relation in (2.5):
$`mF(m\mathrm{id})(F1)`$ $`=`$ $`m(m\mathrm{id})[(\mathrm{\Delta }\mathrm{id})(F)](F1)`$
$`=`$ $`m(\mathrm{id}m)[(\mathrm{id}\mathrm{\Delta })(F)](1F)`$
$`=`$ $`mF(\mathrm{id}m)(1F)`$
as functions from $`A[[t]]_{[[t]]}A[[t]]_{[[t]]}A[[t]]`$ to $`A[[t]]`$. Note that the first relation in (2.4) and the second relation in (2.5) imply that $`1_A`$ remains the multiplicative identity under $``$. ∎
We will need the following notation. Let $`q^\times `$. For every integer $`i1`$, let $`(i)_q=1+q+q^2+\mathrm{}q^{i1},`$ and set $`(0)_q=0`$. Let $`(i)_q!=(i)_q(i1)_q\mathrm{}(1)_q`$ and $`(0)_q!=1`$. The $`q`$-binomial coefficients are
$$\left(\genfrac{}{}{0pt}{}{k}{i}\right)_q=\frac{(k)_q!}{(i)_q!(ki)_q!}$$
for any two integers $`ki0`$. The well-known $`q`$-binomial formula states that
(2.7)
$$(y+z)^k=\underset{i=0}{\overset{k}{}}\left(\genfrac{}{}{0pt}{}{k}{i}\right)_qy^iz^{ki}$$
in any $``$-algebra in which $`y,z`$ are elements such that $`zy=qyz`$.
Let $`S=[x_1,\mathrm{},x_n]`$ and $`q^\times `$. As in \[15, Exer. IV.9.4\], define the linear maps $`q`$-differentiation $`_{i,q}:SS`$ by
(2.8)
$$_{i,q}(x_1^{k_1}\mathrm{}x_n^{k_n})=(k_i)_qx_1^{k_1}\mathrm{}x_{i1}^{k_{i1}}x_i^{k_i1}x_{i+1}^{k_{i+1}}\mathrm{}x_n^{k_n}.$$
If $`q=1`$, these are the usual partial differentiation operators. If $`q`$ is an $`\mathrm{}`$th root of unity, then $`_{i,q}^{\mathrm{}}=0`$ as $`(k)_q=0`$ whenever $`k`$ is a multiple of $`\mathrm{}`$. In general, $`_{i,q}`$ is a skew derivation on $`S`$, specifically
$$_{i,q}(rs)=_{i,q}(r)\tau _{i,q}(s)+r_{i,q}(s)$$
for all $`r,sS`$, where $`\tau _{i,q}`$ is the automorphism of $`S`$ defined by $`\tau _{i,q}(x_1^{k_1}\mathrm{}x_n^{k_n})=q^{k_i}x_1^{k_1}\mathrm{}x_n^{k_n}`$. Under some conditions, these skew derivations may be extended to a crossed product of $`S`$ with a group of linear automorphisms, as we will see.
## 3. A universal deformation formula
Let $`q^\times `$ and let $`H`$ be the algebra generated by $`D_1,D_2,\sigma ^{\pm 1}`$, subject to the relations
$$D_1D_2=D_2D_1,q\sigma D_i=D_i\sigma (i=1,2),\sigma \sigma ^1=1=\sigma ^1\sigma .$$
It is straightforward to check that $`H`$ is a Hopf algebra with
$$\begin{array}{ccccccccc}\hfill \mathrm{\Delta }(D_1)& =& D_1\sigma +1D_1,\hfill & \hfill \epsilon (D_1)& =& 0,\hfill & \hfill S(D_1)& =& D_1\sigma ^1,\hfill \\ \hfill \mathrm{\Delta }(D_2)& =& D_21+\sigma D_2,\hfill & \hfill \epsilon (D_2)& =& 0,\hfill & \hfill S(D_2)& =& \sigma ^1D_2,\hfill \\ \hfill \mathrm{\Delta }(\sigma )& =& \sigma \sigma ,\hfill & \hfill \epsilon (\sigma )& =& 1,\hfill & \hfill S(\sigma )& =& \sigma ^1.\hfill \end{array}$$
If $`q`$ is a primitive $`\mathrm{}`$th root of unity ($`\mathrm{}2`$), the ideal $`I`$ generated by $`D_1^{\mathrm{}}`$ and $`D_2^{\mathrm{}}`$ is a Hopf ideal, that is $`\mathrm{\Delta }(I)IH+HI`$, $`\epsilon (I)=0`$ and $`S(I)I`$. Checking the condition on $`\mathrm{\Delta }`$ involves the $`q`$-binomial formula (2.7) and the observation that $`\left({\displaystyle \genfrac{}{}{0pt}{}{\mathrm{}}{i}}\right)_q=0`$ whenever $`1i\mathrm{}1`$. Thus the quotient $`H/I`$ is also a Hopf algebra. Let
(3.1)
$$H_q=\{\begin{array}{cc}\hfill H/I,& \text{ if }q\text{ is a primitive }\mathrm{}\text{th root of unity }(\mathrm{}2)\hfill \\ \hfill H,& \text{ if }q=1\text{ or is not a root of unity.}\hfill \end{array}$$
We will need the following lemma to obtain a universal deformation formula based on $`H_q`$. If $`q=1`$ or is not a root of unity, we define the $`q`$-exponential function by
$$\mathrm{exp}_q(y)=\underset{i=0}{\overset{\mathrm{}}{}}\frac{1}{(i)_q!}y^i$$
for any element $`y`$ of a $``$-algebra in which this sum is defined. In the proof of Theorem 3.3 below, the $``$-algebra will be $`(H_qH_qH_q)[[t]]`$. If $`q1`$ is a root of unity, this formula makes no sense as some denominators will be zero. We modify the formula as follows in this case. Suppose $`q`$ is a primitive $`\mathrm{}`$th root of unity for $`\mathrm{}2`$. Then we define
$$\mathrm{exp}_q(y)=\underset{i=0}{\overset{\mathrm{}1}{}}\frac{1}{(i)_q!}y^i$$
for any element $`y`$ of a $``$-algebra.
###### Lemma 3.2.
Suppose $`\mathrm{}2`$, $`q`$ is a primitive $`\mathrm{}`$th root of unity, and $`y,z`$ are elements of a $``$-algebra such that $`zy=qyz`$ and $`y^iz^\mathrm{}i=0`$ for $`0i\mathrm{}`$. Then
$$\mathrm{exp}_q(y+z)=\mathrm{exp}_q(y)\mathrm{exp}_q(z).$$
###### Proof.
By the assumed relations and the $`q`$-binomial formula (2.7), each side of the desired equation may be written in the form
$$\underset{i,j=0}{\overset{\mathrm{}1}{}}\frac{1}{(i)_q!(j)_q!}y^iz^j,$$
and thus they are equal. ∎
We note that if $`q`$ is not a primitive root of 1, it is a standard result that $`\mathrm{exp}_q(y+z)=\mathrm{exp}_q(y)\mathrm{exp}_q(z)`$ whenever $`zy=qyz`$ and the relevant sums are defined (see for example \[15, Prop. IV.2.4\]). In the root of unity case, the additional hypothesis stated in the above lemma is required.
###### Theorem 3.3.
Let $`q^\times `$ and let $`H_q`$ be the Hopf algebra defined in (3.1). Then $`\mathrm{exp}_q(tD_1D_2)`$ is a universal deformation formula based on $`H_q`$.
###### Proof.
In case $`q`$ is not a root of unity, this is an unpublished result of Giaquinto and Zhang . Their proof may be adapted to the case $`q`$ is a root of unity by using Lemma 3.2 as follows. (The proof in case $`q`$ is not a root of unity is essentially the same.) Note that the hypotheses of Lemma 3.2 hold for the pairs $`y=tD_1\sigma D_2`$, $`z=tD_1D_2`$ and $`y=tD_1\sigma D_2`$, $`z=tD_1D_21`$ as $`D_2^{\mathrm{}}=0`$ and $`D_1^{\mathrm{}}=0`$, respectively. As $`\mathrm{\Delta }`$ is an algebra homomorphism and $`D_1`$ commutes with $`D_2`$, we thus have
$`(\mathrm{\Delta }\mathrm{id})(\mathrm{exp}_q(tD_1D_2))[\mathrm{exp}_q(tD_1D_2)1]`$
$`=`$ $`\mathrm{exp}_q((\mathrm{\Delta }\mathrm{id})(tD_1D_2))\mathrm{exp}_q(tD_1D_21)`$
$`=`$ $`\mathrm{exp}_q(tD_1\sigma D_2+tD_1D_2)\mathrm{exp}_q(tD_1D_21)`$
$`=`$ $`\mathrm{exp}_q(tD_1\sigma D_2)\mathrm{exp}_q(tD_1D_2)\mathrm{exp}_q(tD_1D_21)`$
$`=`$ $`\mathrm{exp}_q(tD_1\sigma D_2)\mathrm{exp}_q(tD_1D_21)\mathrm{exp}_q(tD_1D_2)`$
$`=`$ $`\mathrm{exp}_q(tD_1\sigma D_2+tD_1D_21)\mathrm{exp}_q(tD_1D_2)`$
$`=`$ $`\mathrm{exp}_q((\mathrm{id}\mathrm{\Delta })(tD_1D_2))(1\mathrm{exp}_q(tD_1D_2))`$
$`=`$ $`(\mathrm{id}\mathrm{\Delta })(\mathrm{exp}_q(tD_1D_2))(1\mathrm{exp}_q(tD_1D_2)).`$
The remaining relation in (2.4) holds as $`\epsilon (D_1)=\epsilon (D_2)=0`$. Thus $`\mathrm{exp}_q(tD_1D_2)`$ is a universal deformation formula. ∎
By Proposition 2.6 and Theorem 3.3, we need only find an $`H_q`$-module algebra $`A`$, and $`m\mathrm{exp}_q(tD_1D_2)`$ will provide a formal deformation of $`A`$. Our first such example is next; a large family of examples is given in Section 4.
###### Example 3.4.
(A Taft algebra.) Let $`A`$ be the algebra defined by generators and relations as follows, where the indices are read modulo $`2`$:
$`A=s_0,s_1,\gamma _0,\gamma _1`$ $``$ $`s_0+s_1=1,s_i^2=s_i,s_is_{i+1}=0,\gamma _i^2=0,\gamma _i\gamma _{i+1}=0,`$
$`s_i\gamma _i=0,s_{i+1}\gamma _i=\gamma _i,\gamma _is_i=\gamma _i,\gamma _is_{i+1}=0.`$
This is an algebra defined by a quiver and relations as in \[5, Thm. 5.1(b)\]; the quiver is Gabriel’s quiver consisting of two arrows in opposite directions between two vertices. The algebra $`A`$ is isomorphic to $`x,ggx=xg,x^2=0,g^2=1`$ via the map $`x\gamma _0\gamma _1`$, $`gs_0s_1`$. It has the structure of a Hopf algebra first discovered by Sweedler, and is one of a series of Hopf algebras constructed by Taft \[19, Example 1.5.6\]. In , Cibils gave deformations of more general classes of algebras defined by quivers and relations, and in this special case one of his deformations may be obtained by applying a universal deformation formula. Specifically, let $`q=1`$ and
$$H_1=D_1,D_2,\sigma ^{\pm 1}D_1D_2=D_2D_1,\sigma D_i=D_i\sigma ,D_i^2=0,\sigma \sigma ^1=1=\sigma ^1\sigma $$
as above. Define
$$\sigma (\gamma _i)=\gamma _{i+1},\sigma (s_i)=s_{i+1},$$
$$D_1(\gamma _i)=s_{i+1},D_2(\gamma _i)=s_i,D_i(s_j)=0.$$
It may be checked that the relations of $`H_1`$ are preserved on the generators of $`A`$, making the vector space $`V=\mathrm{Span}_{}\{s_0,s_1,\gamma _0,\gamma _1\}`$ into an $`H_1`$-module. Therefore the tensor algebra $`T(V)`$ is an $`H_1`$-module algebra, where the action of $`H_1`$ is extended to $`T(V)`$ by (2.5). As $`A`$ is a quotient of $`T(V)`$, it remains to check that the relations of $`A`$ are preserved by the generators of $`H`$, a straightforward computation. (In fact, $`A`$ is also a module algebra for the finite dimensional quotient $`H_1/(\sigma ^21)`$.) By Proposition 2.6 and Theorem 3.3, $`\mathrm{exp}_1(tD_1D_2)=1+tD_1D_2`$ yields a formal deformation of $`A`$. The deformation is
$`A_t=s_0,s_1,\gamma _0,\gamma _1`$ $``$ $`s_0+s_1=1,s_i^2=s_i,s_is_{i+1}=0,\gamma _i^2=0,s_i\gamma _i=0,`$
$`s_{i+1}\gamma _i=\gamma _i,\gamma _is_i=\gamma _i,\gamma _is_{i+1}=0,\gamma _i\gamma _{i+1}=ts_{i+1},`$
which is precisely that given in \[5, Thm. 5.1(b)\]. This deformation is nontrivial since if we specialize to $`t0`$, $`A_t`$ is isomorphic to the $`2\times 2`$ matrix algebra, and thus is not isomorphic to $`A`$. We do not know whether any of Cibils’ other deformations are given by universal deformation formulas.
More generally, suppose that $`A`$ is any $`H_q`$-module algebra. Due to (2.5) and the nature of the coproducts of $`D_1,D_2`$, the following general lemma implies that $`\mu _1=m(D_1D_2)`$ is a Hochschild two-cocycle on $`A`$, that is it satisfies (2.3). This generalizes the well-known fact that the cup product of derivations is a Hochschild two-cocycle. The lemma is proved by direct computation, with no assumption made on the relations among $`D_1,D_2,\sigma `$. If the relations of $`H_q`$ do hold however, then Theorem 3.3 gives an alternative proof that $`\mu _1=m(D_1D_2)`$ is a Hochschild two-cocycle.
###### Lemma 3.5.
Let $`A`$ be an algebra over a field, with multiplication $`m:AAA`$. Let $`\sigma `$ be an automorphism of $`A`$, $`D_1`$ a $`\sigma ,1`$-skew derivation and $`D_2`$ a $`1,\sigma `$-skew derivation of $`A`$. Then $`\mu _1=m(D_1D_2)`$ is a Hochschild two-cocycle.
We end this section with a construction due to Kharchenko \[16, §6.5.5\] of a Hopf algebra of automorphisms and skew derivations of an algebra $`A`$. The Hopf algebras $`H_q`$ are related to some of Kharchenko’s Hopf algebras, and it may be useful to consider his general construction in questions regarding deformations of algebras.
Let $`K`$ be a subgroup of $`\mathrm{Aut}_{}(A)`$. For each $`kK`$, let $`L_k`$ be a vector subspace of $`\mathrm{End}_{}(A)`$ consisting of $`1,k`$-skew derivations of $`A`$. Let $`L=_{kK}L_k`$, and assume $`K`$ acts on $`L`$ in such a way that $`m(L_k)=L_{mkm^1}`$ for all $`k,mK`$. Thus $`K`$ acts by automorphisms on the tensor algebra $`T(L)`$, and we let $`H=T(L)\mathrm{\#}K`$. The coproducts
$$\mathrm{\Delta }(k)=kk\text{ and }\mathrm{\Delta }(D)=D1+kD$$
for all $`kK`$ and $`DL_k`$ extend, by requiring $`\mathrm{\Delta }`$ to be an algebra homomorphism, to a coproduct $`\mathrm{\Delta }`$ on $`H`$. Similarly, the counit $`\epsilon `$ and antipode $`S`$ defined as follows on generators extend to $`H`$: $`\epsilon (k)=1`$, $`\epsilon (D)=0`$, $`S(k)=k^1`$ and $`S(D)=k^1D`$ for all $`kK`$ and $`DL_k`$. Thus $`H`$ is a Hopf algebra.
We obtain the Hopf algebras $`H_q`$ by this construction in the following way: If $`A`$ is an $`H_q`$-module algebra, let $`K`$ be the group generated by the action of $`\sigma `$ on $`A`$, $`L_\sigma =\mathrm{Span}_{}\{D_2\}`$, $`L_{\sigma ^1}=\mathrm{Span}_{}\{D_1\sigma ^1\}`$, and $`L_\tau =\{0\}`$ if $`t\sigma ^{\pm 1}`$. If $`\sigma `$ has infinite order as an automorphism of $`A`$, then $`H_q`$ is a quotient of Kharchenko’s Hopf algebra defined by this data. Otherwise we must take a quotient of $`H_q`$, in which $`\sigma `$ has the correct order, to obtain a quotient of Kharchenko’s Hopf algebra. See for further details on this construction.
## 4. Deformations of group crossed products
In this section we give a large family of group crossed products to which the formula of Theorem 3.3 applies to yield nontrivial deformations.
Let $`G`$ be a group with a representation on a $``$-vector space $`V`$ of dimension $`n`$, so that $`G`$ acts by automorphisms on the symmetric algebra $`S(V)`$. We will identify $`S(V)`$ with polynomials in the variables $`x_1,\mathrm{},x_n`$. In this section, we will be interested in formal deformations of a crossed product $`S(V)\mathrm{\#}_\alpha G`$ for which the infinitesimal $`\mu _1`$ satisfies $`\mu _1(VV)S(V)\overline{g}`$ for some $`gG`$. Not all elements $`gG`$ correspond to such noncoboundary infinitesimals $`\mu _1`$. In case $`G`$ is finite, examination of Hochschild cohomology (see Corollary 6.5 and subsequent comments) shows that we may assume such an element $`g`$ has determinant 1 as an operator on $`V`$, and $`\mathrm{codim}(V^g)=0`$ or $`2`$, where $`V^g=\{vVg(v)=v\}`$, the subspace of $`V`$ invariant under $`g`$. In this section, we will make this assumption, and in addition will assume that $`g`$ is central in $`G`$. Again if the order of $`g`$ is finite, $`g`$ acts diagonally with respect to some basis of $`V`$, and without loss of generality this is $`x_1,\mathrm{},x_n`$. Specifically, we will assume that
(4.1)
$$g(x_1)=qx_1,g(x_2)=q^1x_2,g(x_3)=x_3,\mathrm{},g(x_n)=x_n$$
for some $`q^\times `$. In order to include some infinite groups, we will not assume that $`q`$ is a root of unity. To obtain explicit formulas, we will further need to make a more restrictive assumption:
(4.2)
$$G\text{ preserves the subspaces }x_1,x_2\text{ of }V.$$
If $`q\pm 1`$, this is automatically the case by the assumed centrality of $`g`$. Under the assumption (4.2), we may abuse notation and define the functions $`x_i:G^\times `$ ($`i=1,2`$) by
$$h(x_i)=x_i(h)x_i$$
for each $`hG`$. Let $`D_1,D_2`$ and $`\sigma :S(V)\mathrm{\#}_\alpha GS(V)\mathrm{\#}_\alpha G`$ be the linear functions defined on a basis $`\{x_1^{k_1}\mathrm{}x_n^{k_n}\overline{h}k_i^0,hG\}`$ as follows:
(4.3) $`D_1(x_1^{k_1}\mathrm{}x_n^{k_n}\overline{h})`$ $`=`$ $`x_1(h^1)_{1,q}(x_1^{k_1}\mathrm{}x_n^{k_n})\overline{h},`$
(4.4) $`D_2(x_1^{k_1}\mathrm{}x_n^{k_n}\overline{h})`$ $`=`$ $`q^{k_1}_{2,q^1}(x_1^{k_1}\mathrm{}x_n^{k_n})s\overline{g}\overline{h},`$
(4.5) $`\sigma (x_1^{k_1}\mathrm{}x_n^{k_n}\overline{h})`$ $`=`$ $`x_1(h^1)q^{k_1}x_1^{k_1}\mathrm{}x_n^{k_n}\overline{h},`$
where $`_{1,q}`$, $`_{2,q^1}`$ are defined in (2.8) and $`s[x_3,\mathrm{},x_n]`$ satisfies
(4.6)
$$h(s)=x_1(h)x_2(h)\alpha (g,h)\alpha ^1(h,g)s$$
for all $`hG`$, that is $`s`$ is a semi-invariant of $`G`$. (Our condition on the polynomial $`s`$ is informed by knowledge of Hochschild cohomology; see Corollary 6.5 and the computation (4.11) below.) Calculations using (4.1)–(4.6) and centrality of $`g`$ show that $`\sigma `$ is an automorphism and $`D_1,D_2`$ are skew derivations with respect to $`\sigma `$, specifically
$$D_1(ab)=D_1(a)\sigma (b)+aD_1(b)\text{ and }D_2(ab)=D_2(a)b+\sigma (a)D_2(b)$$
for all $`a,bS(V)\mathrm{\#}_\alpha G`$. A direct calculation shows that $`\mu _1=m(D_1D_2)`$ is a Hochschild two-cocycle on $`S(V)\mathrm{\#}_\alpha G`$, that is $`\mu _1`$ satisfies (2.3). This is also a consequence of Lemma 3.5, or of Theorem 3.3 in combination with Theorem 4.7 below.
Taking $`D_1,D_2`$ to be the skew derivations defined in (4.3), (4.4), the corresponding Hochschild two-cocycle $`\mu _1`$ takes $`VV`$ to $`S(V)\overline{g}`$, the $`g`$-component of $`S(V)\mathrm{\#}_\alpha G`$. If $`g`$ were not central in $`G`$, an associated Hochschild two-cocycle would necessarily involve all components of $`S(V)\mathrm{\#}_\alpha G`$ corresponding to the elements of the conjugacy class of $`g`$ (see Corollary 6.5). We do not know if the explicit formulas of this section can be generalized to noncentral $`g`$.
Let $`H_q`$ be the Hopf algebra defined in (3.1).
###### Theorem 4.7.
Let $`g`$ be a central element of $`G`$ such that (4.1) and (4.2) hold. Then $`S(V)\mathrm{\#}_\alpha G`$ is an $`H_q`$-module algebra under the action defined in (4.3)–(4.5).
###### Proof.
The relations among the generators in $`H_q`$ may be checked to be preserved under the action, so that $`S(V)\mathrm{\#}_\alpha G`$ is an $`H_q`$-module. In particular, in case $`q`$ is a primitive $`\mathrm{}`$th root of unity, $`D_1^{\mathrm{}}=0=D_2^{\mathrm{}}`$ as $`(k)_q=0`$ whenever $`k`$ is a multiple of $`\mathrm{}`$. As stated earlier, $`\sigma `$ is an automorphism of $`S(V)`$ and $`D_1`$ and $`D_2`$ are skew derivations. Clearly $`D_1(1)=0=D_2(1)`$ as $`(0)_q=0`$. Therefore (2.5) holds, so $`S(V)\mathrm{\#}_\alpha G`$ is an $`H_q`$-module algebra. ∎
Combining Proposition 2.6 and Theorems 3.3 and 4.7, we now have the following corollary. In case $`G`$ is finite, the deformations in the corollary are shown to be nontrivial in the remainder of this section. We expect that the same is also true in case $`G`$ is infinite.
###### Corollary 4.8.
Let $`g`$ be a central element of $`G`$ such that (4.1) and (4.2) hold. Then $`\mathrm{exp}_q(tD_1D_2)`$ yields a formal deformation of $`S(V)\mathrm{\#}_\alpha G`$.
We point out that if $`g=1`$ then $`\mathrm{exp}_q(tD_1D_2)`$ restricts to a classical formula on $`S(V)`$, namely
$$\underset{i=0}{\overset{\mathrm{}}{}}\frac{t^i}{i!}\left(\frac{}{x_1}\right)^i\left(s\frac{}{x_2}\right)^i,$$
where $`s[x_3,\mathrm{},x_n]`$ satisfies $`h(s)=x_1(h)x_2(h)s`$ for all $`hG`$. Taking $`G`$ to be the identity group, $`n=2`$, and $`s=1`$, this formula applied to $`[x_1,x_2]`$ yields the Weyl algebra on two generators.
In the special case $`G=/2\times /2`$ and $`V=^3`$ with a particular diagonal action of $`G`$, the formula and deformation of Corollary 4.8 were obtained in \[4, §6\]. The deformations in that case are nontrivial since their corresponding Hochschild two-cocycles are not coboundaries, a consequence of the computations in . Similarly, we now show that the same is true in the more general setting of a finite group $`G`$, based on a computation of the Hochschild cohomology of $`S(V)\mathrm{\#}_\alpha G`$. The Hochschild cohomology was computed by Farinati, Ginzburg and Kaledin in the case $`\alpha =1`$ . The addition of a nontrivial cocycle $`\alpha `$ poses no difficulties, however we need to use some of the details from an explicit algebraic computation. These we provide in the appendix.
There is a chain map from the bar complex (6.1) for $`A=S(V)`$ to the Koszul complex $`K(\{x_i11x_i\}_{i=1}^n)`$,
$$\begin{array}{ccccccccccc}\mathrm{}& & S(V)^4& \stackrel{\delta _2}{}& S(V)^3& \stackrel{\delta _1}{}& S(V)^e& \stackrel{m}{}& S(V)& & 0\\ & & \psi _2& & \psi _1& & & & & & \\ \mathrm{}& & ^2(V)S(V)^e& \stackrel{d_2}{}& ^1(V)S(V)^e& \stackrel{d_1}{}& S(V)^e& \stackrel{m}{}& S(V)& & 0\end{array}$$
We will need an explicit formula for $`\psi _2`$ in particular. A straightforward computation shows that the following formulas work (cf. , in which slightly different formulas are given in the case $`n=3`$):
(4.9)
$$\psi _1(1x_1^{k_1}\mathrm{}x_n^{k_n}1)=\underset{i=1}{\overset{n}{}}\underset{p=1}{\overset{k_i}{}}e_ix_i^{k_ip}x_{i+1}^{k_{i+1}}\mathrm{}x_n^{k_n}x_1^{k_1}\mathrm{}x_{i1}^{k_{i1}}x_i^{p1},$$
(4.10)
$$\psi _2(1x_1^{k_1}\mathrm{}x_n^{k_n}x_1^{m_1}\mathrm{}x_n^{m_n}1)=$$
$$\underset{1i<jn}{}\underset{r=1}{\overset{m_j}{}}\underset{p=1}{\overset{k_i}{}}e_ie_jx_i^{k_ip}x_{i+1}^{k_{i+1}}\mathrm{}x_{j1}^{k_{j1}}x_j^{k_j+m_jr}x_{j+1}^{k_{j+1}+m_{j+1}}\mathrm{}x_n^{k_n+m_n}$$
$$x_1^{k_1+m_1}\mathrm{}x_{i1}^{k_{i1}+m_{i1}}x_i^{m_i+p1}x_{i+1}^{m_{i+1}}\mathrm{}x_{j1}^{m_{j1}}x_j^{r1}.$$
Now assume $`G`$ is a finite group acting on $`V`$, and $`g`$ is a central element of $`G`$ satisfying (4.1) and (4.2) where $`q`$ is a primitive $`\mathrm{}`$th root of unity, $`\mathrm{}2`$. Under these assumptions, by Proposition 6.4 and Corollary 6.5, $`\mathrm{HH}^2(S(V)\mathrm{\#}_\alpha G)`$ contains as the $`g`$-component
$$\mathrm{HH}^2(S(V),S(V)\overline{g})^G(det(\mathrm{Span}_{}\{x_1,x_2\}^{})[x_3,\mathrm{}x_n]\overline{g})^G.$$
Let $`s[x_3,\mathrm{},x_n]\{0\}`$ satisfy (4.6), that is $`h(s)=x_1(h)x_2(h)\alpha (g,h)\alpha ^1(h,g)s`$ for all $`hG`$. Identify the dual function $`(e_1e_2)^{}`$ with a basis of the one-dimensional space $`det(\mathrm{Span}_{}\{x_1,x_2\}^{})`$, where the notation $`e_i`$ comes from the Koszul complex and is defined in the appendix. We first claim that $`(e_1e_2)^{}s\overline{g}`$ corresponds to a nonzero element of $`\mathrm{HH}^2(S(V)\mathrm{\#}_\alpha G)`$ under the above isomorphism. We need only show that $`(e_1e_2)^{}s\overline{g}`$ is invariant under the action of $`G`$. Let $`hG`$. Then
(4.11)
$$h((e_1e_2)^{}s\overline{g})=x_1(h^1)x_2(h^1)(e_1e_2)^{}(h(s))\overline{h}\overline{g}(\overline{h})^1$$
$`=`$ $`\alpha (g,h)\alpha ^1(h,g)(e_1e_2)^{}s\overline{h}\overline{g}(\overline{h})^1`$
$`=`$ $`\alpha (g,h)(e_1e_2)^{}s\overline{hg}(\overline{h})^1`$
$`=`$ $`\alpha (g,h)\alpha ^1(h,h^1)\alpha (hg,h^1)(e_1e_2)^{}s\overline{hgh^1}`$
$`=`$ $`(e_1e_2)^{}s\overline{g}`$
by an application of the two-cocycle identity (2.1) to the triple $`g,h,h^1`$, since $`gC(G)`$.
Next we show that the nonzero element $`(e_1e_2)^{}s\overline{g}`$ of $`\mathrm{HH}^2(S(V)\mathrm{\#}_\alpha G)`$ may be identified with a Hochschild two-cocycle $`\mu _1`$ of the form $`m(D_1D_2)`$ where $`D_1,D_2`$ are defined in (4.3), (4.4). This will prove that the formal deformations of $`S(V)\mathrm{\#}_\alpha G`$ given by Corollary 4.8 are nontrivial in case $`G`$ is finite. We will need \[4, Thm. 5.4\], which will be applied to a Koszul resolution:
###### Proposition 4.12 (Caldararu-Giaquinto-Witherspoon).
Let $`A=S(V)\mathrm{\#}_\alpha G`$. Let $`f:P_nA`$ be a function representing an element of $`\mathrm{HH}^n(S(V),A)^G\mathrm{HH}^n(A)`$ expressed in terms of any $`S(V)^e`$-projective resolution $`P_{\text{}}`$ of $`S(V)`$ carrying an action of $`G`$. The corresponding function $`\stackrel{~}{f}\mathrm{Hom}_{}(A^n,A)\mathrm{Hom}_{A^e}(A^{(n+2)},A)`$ on the bar complex (6.1) is given by
$$\stackrel{~}{f}(p_1\overline{\sigma _1}\mathrm{}p_n\overline{\sigma _n})=((f\psi _n)(1p_1\sigma _1(p_2)\mathrm{}(\sigma _1\mathrm{}\sigma _{n1})(p_n)1))\overline{\sigma _1}\mathrm{}\overline{\sigma _n}.$$
In particular, if $`n=2`$, we obtain the infinitesimal deformation $`\mu _1:AAA`$,
$$\mu _1(p_1\overline{\sigma _1}p_2\overline{\sigma _2})=((f\psi _2)(1p_1\sigma _1(p_2)1))\overline{\sigma _1}\overline{\sigma _2}.$$
As a consequence of the proposition, the element $`(e_1e_2)^{}s\overline{g}`$ of $`\mathrm{HH}^2(S(V)\mathrm{\#}_\alpha G)`$ may be identified with the function $`\mu _1:AAA`$ where $`\mu _1(x_1^{k_1}\mathrm{}x_n^{k_n}\overline{h}x_1^{m_1}\mathrm{}x_n^{m_n}\overline{k})`$ is $`\psi _2(1x_1^{k_1}\mathrm{}x_n^{k_n}h(x_1^{m_1}\mathrm{}x_n^{m_n})1)`$ followed by application of the function representing $`(e_1e_2)^{}s\overline{g}`$ at the chain level, and right multiplication by $`\overline{h}\overline{k}`$. By our hypotheses, we have
$`\psi _2(1x_1^{k_1}\mathrm{}x_n^{k_n}h(x_1^{m_1}\mathrm{}x_n^{m_n})1)`$
$$=x_1(h)^{m_1}x_2(h)^{m_2}\psi _2(1x_1^{k_1}\mathrm{}x_n^{k_n}x_1^{m_1}x_2^{m_2}h(x_3^{m_3}\mathrm{}x_n^{m_n})).$$
By (4.10), the resulting coefficient of $`e_1e_2`$ is
$$x_1(h)^{m_1}x_2(h)^{m_2}\underset{r=1}{\overset{m_2}{}}\underset{p=1}{\overset{k_1}{}}x_1^{k_1p}x_2^{k_2+m_2r}x_3^{k_3}\mathrm{}x_n^{k_n}h(x_3^{m_3}\mathrm{}x_n^{m_n})x_1^{m_1+p1}x_2^{r1}.$$
Applying $`(e_1e_2)^{}s\overline{g}`$ and multiplying by $`\overline{h}\overline{k}`$, we obtain
$`x_1(h)^{m_1}x_2(h)^{m_2}{\displaystyle \underset{r=1}{\overset{m_2}{}}}{\displaystyle \underset{p=1}{\overset{k_1}{}}}x_1^{k_1p}x_2^{k_2+m_2r}x_3^{k_3}\mathrm{}x_n^{k_n}h(x_3^{m_3}\mathrm{}x_n^{m_n})s\overline{g}x_1^{m_1+p1}x_2^{r1}\overline{h}\overline{k}`$
$`=`$ $`x_1(h)^{m_1}x_2(h)^{m_2}{\displaystyle \underset{r=1}{\overset{m_2}{}}}{\displaystyle \underset{p=1}{\overset{k_1}{}}}q^{m_1+pr}x_1^{k_1+m_11}x_2^{k_2+m_21}x_3^{k_3}\mathrm{}x_n^{k_n}h(x_3^{m_3}\mathrm{}x_n^{m_n})s\overline{g}\overline{h}\overline{k}`$
$`=`$ $`q^{m_1}x_1(h)^{m_1}x_2(h)^{m_2}(k_1)_q(m_2)_{q^1}x_1^{k_1+m_11}x_2^{k_2+m_21}x_3^{k_3}\mathrm{}x_n^{k_n}h(x_3^{m_3}\mathrm{}x_n^{m_n})s\overline{g}\overline{h}\overline{k}.`$
On the other hand,
$`(m(D_1D_2))(x_1^{k_1}\mathrm{}x_n^{k_n}\overline{h}x_1^{m_1}\mathrm{}x_n^{m_n}\overline{k})`$
$`=`$ $`x_1(h^1)(k_1)_qx_1^{k_11}x_2^{k_2}\mathrm{}x_n^{k_n}\overline{h}q^{m_1}(m_2)_{q^1}x_1^{m_1}x_2^{m_21}x_3^{m_3}\mathrm{}x_n^{m_n}s\overline{g}\overline{k}`$
$`=`$ $`q^{m_1}x_1(h)^{m_1}x_2(h)^{m_2}(k_1)_q(m_2)_{q^1}x_1^{k_1+m_11}x_2^{k_2+m_21}x_3^{k_3}\mathrm{}x_n^{k_n}h(x_3^{m_3}\mathrm{}x_n^{m_n})s\overline{g}\overline{h}\overline{k}.`$
Therefore $`m(D_1D_2)`$ is the Hochschild two-cocycle represented by $`(e_1e_2)^{}s\overline{g}`$ and so is not a coboundary. This implies that the formal deformations given by Corollary 4.8 are nontrivial in case $`G`$ is finite.
###### Example 4.13.
(Twisted graded Hecke algebras.) Let $`G`$ be a finite subgroup of $`GL(V)`$, and $`gG`$ a central element satisfying (4.1) and (4.2). Suppose $`s=1`$ satisfies (4.6), that is $`1=x_1(h)x_2(h)\alpha (g,h)\alpha ^1(h,g)`$ for all $`hG`$. We may rewrite this condition as $`det(h|_{(V^g)^{}})=\alpha (h,g)\alpha ^1(g,h)`$. Then $`\mu _1=m(D_1D_2)`$ is a bilinear form on $`S(V)\mathrm{\#}_\alpha G`$ of degree $`2`$, where $`S(V)\mathrm{\#}_\alpha G`$ is a graded algebra in which elements of $`V`$ have degree 1 and elements of $`G`$ have degree 0. More generally, in the formula $`\mathrm{exp}_q(tD_1D_2)`$, the bilinear form $`\mu _i=\frac{1}{(i)_q!}m(D_1^iD_2^i)`$ has degree $`2i`$. By \[25, Thm 3.2\], the resulting formal deformation of $`S(V)\mathrm{\#}_\alpha G`$ becomes a (twisted) graded Hecke algebra when the scalars are restricted to $`[t]`$. In this case, that means the associated deformation of $`S(V)\mathrm{\#}_\alpha G`$ over $`[t]`$ is isomorphic to
$$T(V)\mathrm{\#}_\alpha G[t]/(vwwva_g(v,w)t\overline{g}),$$
the quotient by the ideal generated by all elements $`vwwva_g(v,w)t\overline{g}`$, for $`v,wV`$, where $`a_g(v,w)=\mu _1(v,w)\mu _1(w,v)`$. This (twisted) graded Hecke algebra is special in that only one such function $`a_g`$ is nonzero. In the next section, we give some examples for which there is an analogous deformation with more than one group element $`g`$ having $`a_g`$ nonzero.
## 5. Universal Deformations
In this section we give examples for which some of the universal deformation formulas from the last section, corresponding to different group elements, may be combined into larger formulas. The first example generalizes \[4, Lemma 6.2\].
###### Example 5.1.
Let $`n3`$ and $`\mathrm{}2`$ be integers, $`q`$ a primitive $`\mathrm{}`$th root of unity, $`G=(/\mathrm{})^{n1}`$ and $`V=^n`$. Identify $`G`$ with the subgroup of $`\mathrm{SL}(V)`$ generated by the diagonal matrices
$`g_1`$ $`=`$ $`\mathrm{diag}(q,q^1,1,\mathrm{},1),`$
$`g_2`$ $`=`$ $`\mathrm{diag}(1,q,q^1,1,\mathrm{},1),`$
$`\mathrm{}`$
$`g_{n1}`$ $`=`$ $`\mathrm{diag}(1,\mathrm{},1,q,q^1),`$
with respect to a basis $`x_1,\mathrm{},x_n`$ of $`V`$. Let $`g_n=g_1^1\mathrm{}g_{n1}^1=\mathrm{diag}(q^1,1,\mathrm{},1,q)`$. Let $`\alpha :G\times G^\times `$ be the following two-cocycle:
$$\alpha (g_1^{i_1}\mathrm{}g_{n1}^{i_{n1}},g_1^{j_1}\mathrm{}g_{n1}^{j_{n1}})=q^{_{1kn2}i_kj_{k+1}}.$$
It may be checked directly that $`\alpha `$ satisfies the two-cocycle condition (2.1). Note that $`\alpha `$ is not a coboundary: By their definition, two-coboundaries for abelian groups are symmetric, but $`\alpha `$ is clearly not symmetric. In case $`n=3,\mathrm{}=2`$, $`\alpha `$ is cohomologous to the nontrivial cocycle given in \[4, Example 3.4\], as in that case there is a unique nontrivial two-cocycle up to coboundary. Note that
(5.2)
$$\alpha (g_{i+1},g_i)=q\alpha (g_i,g_{i+1})\text{for }1in$$
(where $`g_{n+1}=g_1`$ by definition). Direct calculations also show that
(5.3)
$$h(x_ix_{i+1})=\frac{\alpha (h,g_i)}{\alpha (g_i,h)}x_ix_{i+1}\text{for }1in,$$
for all $`hG`$ (where $`x_{n+1}=x_1`$ by definition).
Let $`H_i=H_q`$ ($`1in`$) be the Hopf algebra defined in (3.1), acting on $`S(V)\mathrm{\#}_\alpha G`$ via the formulas (4.3)–(4.5), where we replace $`g`$ by $`g_i`$ and $`x_1,x_2`$ by $`x_i,x_{i+1}`$. Applying (4.6) and (5.3), it may be checked that the polynomial $`s`$ arising in the action of $`H_i`$ must be in $`S(V^{g_i})(S(V))^G=[x_1^{\mathrm{}},\mathrm{},x_{i1}^{\mathrm{}},x_{i+2}^{\mathrm{}},\mathrm{},x_n^{\mathrm{}}]`$, where if $`i=n`$ we leave out $`x_1`$ and $`x_n`$. Using this fact, equation (5.2), and the identity $`(i+\mathrm{})_{q^1}=(i)_{q^1}`$ for all integers $`i`$, it may be checked directly that the images of the $`H_i`$ ($`1in`$) in $`\mathrm{End}_{}(S(V)\mathrm{\#}_\alpha G)`$ mutually commute. Thus there is a corresponding algebra homomorphism from the Hopf algebra $`H_1\mathrm{}H_n`$ to $`\mathrm{End}_{}(S(V)\mathrm{\#}_\alpha G)`$, and $`S(V)\mathrm{\#}_\alpha G`$ is a module algebra for $`H_1\mathrm{}H_n`$. A product of universal deformation formulas is again a universal deformation formula, based on the tensor product of the bialgebras. Thus by Theorem 3.3,
(5.4)
$$\mathrm{exp}_q(tD_1^{g_1}D_2^{g_1})\mathrm{}\mathrm{exp}_q(tD_1^{g_n}D_2^{g_n})$$
(where superscripts indicate the Hopf subalgebra from which the operators originate) is a universal deformation formula based on $`H_1\mathrm{}H_n`$. By Proposition 2.6, this formula applies to yield a formal deformation of $`S(V)\mathrm{\#}_\alpha G`$.
In case $`s=1`$, restricting this deformation to one over $`[t]`$ results in the twisted graded Hecke algebra
$$T(V)\mathrm{\#}_\alpha G[t]/\left(vwwv\underset{i=1}{\overset{n}{}}a_{g_i}(v,w)t\overline{g}_i\right)$$
where $`a_{g_i}(v,w)=D_1^{g_i}(v)D_2^{g_i}(w)D_1^{g_i}(w)D_2^{g_i}(v)`$. The scalar coefficients of the $`a_{g_i}`$ may be varied independently to obtain a vector space of dimension $`n`$ parametrizing the possible twisted graded Hecke algebras realizable by the formula (5.4) and scalar modifications. It is shown in \[25, Example 2.16\] that these are in fact all the twisted graded Hecke algebras for this choice of $`G`$ and $`\alpha `$ in case $`\mathrm{}2`$.
In case $`\mathrm{}=2`$ and $`n=3`$, the elements $`g_1,g_2,g_3`$ are precisely the nonidentity elements of $`G`$. The formal deformation of $`S(V)\mathrm{\#}_\alpha G`$ arising from the formula (5.4) is nearly the universal deformation, as is justified by considering the Hochschild cohomology of $`S(V)\mathrm{\#}_\alpha G`$ (see \[4, Example 4.7\] or the more general Corollary 6.5). That is, every Hochschild two-cocycle $`\mu _1`$ with image in $`S(V)\mathrm{\#}_\alpha (G\{1\})`$ is an infinitesimal of the formal deformation resulting from (5.4) with appropriate choices of the polynomials $`s`$ in (4.4). (Classical deformations corresponding to the choice $`g=1`$ involve derivations that do not commute with the actions of the $`H_i`$, and so we do not include these in the formula.) If $`n>3`$ or $`\mathrm{}>2`$, there are nonidentity group elements other than $`g_1,\mathrm{},g_n`$, and the actions of the corresponding Hopf algebras may no longer commute (but see the next example below). If $`\alpha `$ is not taken to be the cocycle we have chosen, the images of the $`H_i`$ in $`\mathrm{End}_{}(S(V)\mathrm{\#}_\alpha G)`$ again may no longer commute, and we do not know whether there is a universal deformation formula more complicated than (5.4) involving these operators.
###### Example 5.5.
Let $`G`$ be a group acting on a vector space $`V`$ of dimension $`n`$, $`g`$ a central element of $`G`$, and assume (4.1) and (4.2) hold with $`q`$ a primitive $`\mathrm{}`$th root of unity, $`\mathrm{}>2`$. Thus $`g`$ corresponds to $`\mathrm{diag}(q,q^1,1,\mathrm{},1)`$ and $`g^1`$ corresponds to $`\mathrm{diag}(q^1,q,1,\mathrm{},1)`$. Assume further that $`\alpha (g,g^1)=\alpha (g^1,g)`$, as is true in the last example for $`g=g_i`$. (In case $`G`$ is finite, this assumption imposes no loss of generality, as any two-cocyle is cohomologous to one satisfying this assumption \[14, Thm. 3.6.2\]). Consider the images of $`H_q`$ and $`H_{q^1}`$ in $`\mathrm{End}_{}(S(V)\mathrm{\#}_\alpha G)`$, where we let $`D_1^{g^1}`$ involve $`q^1`$-differentiation with respect to $`x_1`$ and $`D_2^{g^1}`$ involve $`q`$-differentiation with respect to $`x_2`$ in (4.3) and (4.4). Multiplying and dividing the left side of the equation below by $`q^{i2}q^{i1}`$ yields the right side:
$$\frac{(i)_q(i1)_{q^1}}{(i)_{q^1}(i1)_q}=q.$$
Using this identity and the assumption $`\alpha (g,g^1)=\alpha (g^1,g)`$, we find that the following relations hold in $`\mathrm{End}_{}(S(V)\mathrm{\#}_\alpha G)`$ among the images of the generators of $`H_q`$ and $`H_{q^1}`$:
$`D_1^{g^1}D_1^g=qD_1^gD_1^{g^1}`$ , $`D_2^{g^1}D_2^g=q^1D_2^gD_2^{g^1},`$
$`D_1^gD_2^{g^1}=D_2^{g^1}D_1^g`$ , $`D_1^{g^1}D_2^g=D_2^gD_1^{g^1},`$
$`\sigma ^g\sigma ^{g^1}`$ $`=`$ $`\sigma ^{g^1}\sigma ^g,`$
$`\sigma ^gD_1^{g^1}=q^1D_1^{g^1}\sigma ^g`$ , $`\sigma ^gD_2^{g^1}=qD_2^{g^1}\sigma ^g,`$
$`\sigma ^{g^1}D_1^g=qD_1^g\sigma ^{g^1}`$ , $`\sigma ^{g^1}D_2^g=q^1D_2^g\sigma ^{g^1}.`$
These relations are preserved by $`\mathrm{\Delta }`$, $`\epsilon `$ and $`S`$, and so the algebra generated by $`H_q`$ and $`H_{q^1}`$, subject to the above relations, is a Hopf algebra. The proof of Theorem 3.3 may be modified to show that $`\mathrm{exp}_q(tD_1^gD_2^g)\mathrm{exp}_{q^1}(tD_1^{g^1}D_2^{g^1})`$ is a universal deformation formula. The key idea is to move factors corresponding to $`g^1`$ past factors corresponding to $`g`$, so that the proof of Theorem 3.3 may be applied separately for each of $`g,g^1`$. The relations above imply that indeed the appropriate factors commute. In case $`G=/3`$, this formula will nearly result in a universal deformation (again having infinitesimal with image in $`S(V)\mathrm{\#}_\alpha (G\{1\})`$).
## 6. Appendix: A computation of Hochschild cohomology
The Hochschild cohomology of a $``$-algebra $`A`$ is $`\mathrm{HH}^{\text{}}(A):=\mathrm{Ext}_{A^e}^{\text{}}(A,A)`$, where $`A^e=AA^{op}`$ acts on $`A`$ by left and right multiplication. More generally, if $`M`$ is an $`A`$-bimodule (equivalently, an $`A^e`$-module), we may define $`\mathrm{HH}^{\text{}}(A,M):=\mathrm{Ext}_{A^e}^{\text{}}(A,M)`$, so that $`\mathrm{HH}^{\text{}}(A)=\mathrm{HH}^{\text{}}(A,A)`$. These Ext groups may be expressed via the $`A^e`$-free resolution of $`A`$:
(6.1)
$$\mathrm{}\stackrel{\delta _3}{}A^4\stackrel{\delta _2}{}A^3\stackrel{\delta _1}{}A^e\stackrel{m}{}A0,$$
where $`m`$ is multiplication and
$$\delta _i(a_0a_1\mathrm{}a_{i+1})=\underset{j=0}{\overset{i}{}}(1)^ja_0\mathrm{}a_ja_{j+1}\mathrm{}a_{i+1}.$$
Applying $`\mathrm{Hom}_{A^e}(,M)`$ and dropping the term $`\mathrm{Hom}_{A^e}(A,M)`$, we obtain
$$0\mathrm{Hom}_{A^e}(A^e,M)\stackrel{\delta _1^{}}{}\mathrm{Hom}_{A^e}(A^3,M)\stackrel{\delta _2^{}}{}\mathrm{Hom}_{A^e}(A^4,M)\stackrel{\delta _3^{}}{}\mathrm{}$$
Then $`\mathrm{HH}^i(A,M)=\mathrm{Ker}(\delta _{i+1}^{})/\mathrm{Im}(\delta _i^{})`$ and $`\mathrm{HH}^{\text{}}(A,M)=_{i0}\mathrm{HH}^i(A,M)`$. Noting that $`\mathrm{Hom}_{A^e}(A^{(i+2)},A)\mathrm{Hom}_{}(A^i,A)`$, a straightforward calculation shows that $`\mathrm{HH}^2(A)`$ may be identified with the space of $``$-linear functions $`\mu _1:AAA`$ satisfying the Hochschild two-cocycle condition (2.3), modulo coboundaries. See for more details on Hochschild cohomology.
Let $`G`$ be a finite subgroup of $`\mathrm{GL}(V)`$. We will compute $`\mathrm{HH}^{\text{}}(S(V)\mathrm{\#}_\alpha G)`$, using techniques similar to those in , where the crossed product was taken with a Weyl algebra rather than a polynomial algebra. We will use a result of Ştefan on Hopf Galois extensions \[22, Cor. 3.4\]. It implies that there is an action of $`G`$ on $`\mathrm{HH}^{\text{}}(S(V),S(V)\mathrm{\#}_\alpha G)`$ for which
(6.2)
$$\mathrm{HH}^{\text{}}(S(V)\mathrm{\#}_\alpha G)\mathrm{HH}^{\text{}}(S(V),S(V)\mathrm{\#}_\alpha G)^G,$$
where the superscript $`G`$ denotes the subspace of $`G`$-invariant elements. (A more explicit proof of this result, useful in this context, is given in \[4, §5\]). A Koszul complex may then be used to compute $`\mathrm{HH}^{\text{}}(S(V),S(V)\mathrm{\#}_\alpha G)`$. This is done in a more general geometric setting by Ginzburg and Kaledin in the case $`G`$ is symplectic and $`\alpha `$ is trivial, although they note that their techniques apply to any finite group $`G`$. An elegant algebraic computation is given by Farinati for an arbitrary finite group $`G`$, and trivial $`\alpha `$. The additive structure of $`\mathrm{HH}^{\text{}}(S(V),S(V)\mathrm{\#}_\alpha G)`$, before taking $`G`$-invariants, is independent of $`\alpha `$ since the $`S(V)`$-bimodule structure of $`S(V)\mathrm{\#}_\alpha G`$ does not involve $`\alpha `$. Thus the techniques of either or apply here. For completeness, we give an explicit algebraic computation whose details are needed in Section 4.
Note that $`S(V)\mathrm{\#}_\alpha G=_{gG}S(V)\overline{g}`$, where $`S(V)\overline{g}=\{s\overline{g}sS(V)\}`$, as an $`S(V)`$-bimodule. Thus there is an additive decomposition of Hochschild cohomology,
(6.3)
$$\mathrm{HH}^{\text{}}(S(V),S(V)\mathrm{\#}_\alpha G)\underset{gG}{}\mathrm{HH}^{\text{}}(S(V),S(V)\overline{g}).$$
We will determine each summand $`\mathrm{HH}^{\text{}}(S(V),S(V)\overline{g})`$, noting again that $`\alpha `$ plays no role here as we need only the $`S(V)`$-module structure of each $`S(V)\overline{g}`$. If $`g=1`$, we have $`\mathrm{HH}^{\text{}}(S(V),S(V)\overline{1})=\mathrm{HH}^{\text{}}(S(V))`$, and the Hochschild-Kostant-Rosenberg Theorem states that
$$\mathrm{HH}^{\text{}}(S(V))_{S(V)}^{\text{}}(S(V)^n)^{\text{}}(V^{})S(V),$$
where $`n=dimV`$. (See \[24, Exer. 9.1.3 and Thm. 9.4.7\].) Letting $`x_1,\mathrm{},x_n`$ be a basis of $`V`$, this may be computed directly from the $`S(V)^e`$-projective Koszul resolution $`K(\{x_i11x_i\}_{i=1}^n)^{\text{}}(V)S(V)^e`$ of $`S(V)`$. (See \[24, §4.5\] for details on Koszul complexes.) The differential $`d_m:^m(V)S(V)^e`$ is given by
$$d_m(e_{i_1}\mathrm{}e_{i_m}11)=\underset{k=1}{\overset{m}{}}(1)^{k+1}e_{i_1}\mathrm{}\widehat{e_{i_k}}\mathrm{}e_{i_m}(x_{i_k}11x_{i_k}),$$
where we use the standard notation $`e_{i_j}`$ for the element $`x_{i_j}`$ in $`^1(V)`$. After application of $`\mathrm{Hom}_{S(V)^e}(,S(V))`$, all chain maps become 0 as $`S(V)`$ is commutative. We will identify $`\mathrm{HH}^{\text{}}(S(V))`$ with $`^{\text{}}(V^{})S(V)`$, as the group action is clear in that notation: It is diagonal on the factors, with the standard actions on $`S(V)`$ and on $`^{\text{}}(V^{})`$. In case of an element $`g`$ not necessarily equal to 1, we have the following.
###### Proposition 6.4.
For each $`gG`$,
$$\mathrm{HH}^{\text{}}(S(V),S(V)\overline{g})^{\text{}\mathrm{codim}V^g}((V^g)^{})S(V^g).$$
In particular, the lowest degree $`j`$ for which $`\mathrm{HH}^j(S(V),S(V)\overline{g})0`$ is $`j=\mathrm{codim}V^g`$.
###### Proof.
Fix $`gG`$. As the order of $`g`$ is finite, we may assume without loss of generality that the action of $`g`$ is diagonal with respect to the basis $`x_1,\mathrm{},x_n`$ of $`V`$. Thus there are scalars $`\lambda _i`$ with $`gx_i=\lambda _ix_i(i=1,\mathrm{},n)`$. We will further assume, for notational convenience, that the basis is ordered so that $`\lambda _i=1`$ for $`i=1,\mathrm{},r`$ and $`\lambda _i1`$ for $`i=r+1,\mathrm{},n`$. (Set $`r=0`$ if $`V^g=0`$ and $`r=n`$ if $`V^g=V`$.) We may also assume that $`\mathrm{Span}_{}\{x_{r+1},\mathrm{},x_n\}=(V^g)^{}`$ where the orthogonal complement is taken with respect to some nondegenerate $`G`$-invariant Hermitian form on $`V`$.
Consider the complex $`\mathrm{Hom}_{S(V)^e}(K(\{x_i11x_i\}_{i=1}^n),S(V)\overline{g})`$, which we may identify with $`\mathrm{Hom}_{S(V)^e}(^{\text{}}(V)S(V)^e,S(V)\overline{g})^{\text{}}(V^{})S(V)\overline{g}`$. Additively, this is the same as $`^{\text{}}(V^{})S(V)`$, but the factor $`\overline{g}`$ affects the differentials, which we will determine next. They are not necessarily all zero (in contrast to the case $`g=1`$). If $`sS(V)`$, we have
$$(x_i11x_i)s\overline{g}=x_is\overline{g}s\overline{g}x_i=(x_igx_i)s\overline{g}.$$
If $`i=1,\mathrm{},r`$, this element is 0. If $`i=r+1,\mathrm{},n`$, the factor $`x_igx_i=(1\lambda _i)x_i`$ is a nonzero scalar multiple of $`x_i`$. Thus $`\mathrm{Hom}_{S(V)^e}(K(\{x_i11x_i\}_{i=1}^n),S(V)\overline{g})`$ is equivalent to the dual Koszul complex $`\overline{K}(0,\mathrm{},0,x_{r+1},\mathrm{},x_n)`$ for $`S(V)`$, where the bar denotes the reverse order. This is the tensor product (over $`S(V)`$) of two complexes for $`S(V)`$: $`\overline{K}(0,\mathrm{},0)`$ and $`\overline{K}(x_{r+1},\mathrm{},x_n)`$. The second complex is exact other than in degree $`nr`$ (as the corresponding Koszul complex is exact other than in degree 0), where it has cohomology $`S(V)/(x_{r+1},\mathrm{},x_n)S(V)S(V^g)`$ (see \[24, Cor. 4.5.4\]). We will identify this with $`det(((V^g)^{})^{})S(V^g)`$, where $`det(((V^g)^{})^{})`$ is the one-dimensional space $`^{\mathrm{codim}V^g}(((V^g)^{})^{})`$, to account for the degree shift and the action of $`G`$. The spectral sequence of the double complex $`\overline{K}(0,\mathrm{},0)_{S(V)}\overline{K}(x_{r+1},\mathrm{},x_n)`$ thus collapses at $`E_2`$ with $`E_2^{pq}=0`$ for $`qnr`$, and $`E_2^{p,nr}=\mathrm{H}^p(\overline{K}(0,\mathrm{},0))_{S(V)}S(V^g)`$ by freeness of the terms of the chain complex over $`S(V)`$. This follows from \[3, Thm. 3.4.2\], which also implies that the cohomology is precisely $`E_2^{p,nr}`$. Now $`\mathrm{H}^p(\overline{K}(0,\mathrm{},0))_{S(V)}^p(S(V)^{dimV^g})^p((V^g)^{})S(V)`$, and as this is tensored with the cohomology of $`\overline{K}(x_{r+1},\mathrm{},x_n)`$, namely $`S(V)/(x_{r+1},\mathrm{},x_n)S(V)S(V^g)`$ in degree $`nr`$, we obtain the stated result. ∎
We will identify the cohomology $`^{\text{}\mathrm{codim}V^g}((V^g)^{})S(V^g)`$ of the theorem with $`^{\text{}\mathrm{codim}V^g}((V^g)^{})det(((V^g)^{})^{})S(V^g)\overline{g}`$. The action of $`G`$ is nontrivial on the one-dimensional factor $`det(((V^g)^{})^{})`$, as may be seen by considering the action of $`G`$ on the corresponding cochain complex, and the action on $`\overline{g}`$ is by conjugation by $`\overline{h}`$ ($`hG`$).
The following corollary is immediate from (6.2), (6.3) and Proposition 6.4, after making the above identifications.
###### Corollary 6.5.
There is an additive decomposition of Hochschild cohomology,
$$\mathrm{HH}^{\text{}}(S(V)\mathrm{\#}_\alpha G)\left(\underset{gG}{}^{\text{}\mathrm{codim}V^g}((V^g)^{})det(((V^g)^{})^{})S(V^g)\overline{g}\right)^G.$$
Compare the above corollary with \[9, Thm. 3.6\], or with the formula just above (6.4) in .
As a consequence of Corollary 6.5, we obtain a necessary condition for there to exist a Hochschild two-cocycle $`\mu _1`$, with image in the $`g`$-component $`S(V)\overline{g}`$, that is not a coboundary. Due to the degree shift $`2\mathrm{codim}V^g`$, such an element $`g`$ must satisfy $`\mathrm{codim}V^g\{0,1,2\}`$. We claim that the determinant of $`g`$ on $`V`$ must be 1. If $`det(g)1`$, the action of $`g`$ itself on the one-dimensional space $`det(((V^g)^{})^{})`$ is nontrivial, whereas its actions on $`^{2\mathrm{codim}V^g}((V^g)^{})`$ and on $`S(V^g)\overline{g}`$ are trivial. Consequently there can be no such $`G`$-invariant elements. Therefore $`det(g)=1`$, which also now implies $`\mathrm{codim}V^g\{0,2\}`$. (See also \[9, Ex. 3.10\].) |
warning/0506/math0506360.html | ar5iv | text | # Grothendieck bialgebras, Partition lattices and symmetric functions in noncommutative variables
## Introduction
Combinatorial Hopf algebras are graded connected Hopf algebras equipped with a multiplicative linear functional $`\zeta :\mathrm{𝕜}`$ called a character (see ). Here we assume that $`\mathrm{𝕜}`$ is a field of characteristic zero. There has been renewed interest in these spaces in recent papers (see for example and the references therein) and one particularly interesting aspect of recent work has been to realize a given combinatorial Hopf algebra as the Grothendieck Hopf algebra of a tower of algebras.
The prototypical example is the Hopf algebra of symmetric functions viewed, via the Frobenius characteristic map, as the Grothendieck Hopf algebras of the modules of all symmetric group algebras $`\mathrm{𝕜}S_n`$ for $`n0`$. The multiplication is given via induction from $`\mathrm{𝕜}S_n\mathrm{𝕜}S_m`$ to $`\mathrm{𝕜}S_{n+m}`$ and the comultiplication is the sum over $`r`$ of the restriction from $`\mathrm{𝕜}S_n`$ to $`\mathrm{𝕜}S_r\mathrm{𝕜}S_{nr}`$. The tensor product of modules defines a third operation on symmetric functions usually referred to as the internal multiplication or the Kronecker product . The Schur symmetric functions are then canonically defined as the Frobenius image of the simple modules.
There are many more examples of this kind of connection (see ). Here we are interested in the bialgebra structure of the symmetric functions in noncommutative variables and the goal of this paper is to realize it as the Grothendieck bialgebra of the modules of the partition lattice algebras.
We denote by $`\mathrm{𝖭𝖢𝖲𝗒𝗆}=_{d0}\mathrm{𝖭𝖢𝖲𝗒𝗆}_d`$ the algebra of symmetric functions in noncommutative variables, the product is induced from the concatenation of words. This is a Hopf algebra equipped with an internal comultiplication. The space $`\mathrm{𝖭𝖢𝖲𝗒𝗆}_d`$ is the subspace of series in the noncommutative variables $`x_1,x_2,\mathrm{}`$ with homogeneous degree $`d`$ that are invariants by any finite permutation of the variables. The algebra structure of $`\mathrm{𝖭𝖢𝖲𝗒𝗆}`$ was first introduced in were it was shown to be a free noncommutative algebra. This algebra was used in to study free powers of noncommutative rings. More recently, a series of new bases was given for this space, lifting some of the classical bases of (commutative) symmetric functions . The Hopf algebra structure was uncovered in along with other fundamental algebraic and geometric structures.
The (external) comultiplication $`\mathrm{\Delta }:\mathrm{𝖭𝖢𝖲𝗒𝗆}_d\mathrm{𝖭𝖢𝖲𝗒𝗆}_k\mathrm{𝖭𝖢𝖲𝗒𝗆}_{dk}`$ is graded and gives rise to a structure of a graded Hopf algebra on $`\mathrm{𝖭𝖢𝖲𝗒𝗆}`$. The algebra $`\mathrm{𝖭𝖢𝖲𝗒𝗆}`$ also has an internal comultiplication $`\mathrm{\Delta }^{}:\mathrm{𝖭𝖢𝖲𝗒𝗆}_d\mathrm{𝖭𝖢𝖲𝗒𝗆}_d\mathrm{𝖭𝖢𝖲𝗒𝗆}_d`$ which is not graded. The algebra $`\mathrm{𝖭𝖢𝖲𝗒𝗆}`$ with the comultiplication $`\mathrm{\Delta }^{}`$ is only a bialgebra (not graded) and is different from the previous graded Hopf structure.
After investigating the Hopf algebra structure of $`\mathrm{𝖭𝖢𝖲𝗒𝗆}`$, it is natural to ask if there exists a tower of algebras $`\{A_n\}_{n0}`$ such that the Hopf algebra $`\mathrm{𝖭𝖢𝖲𝗒𝗆}`$ corresponds to the Grothendieck bialgebra (or Hopf) algebra of the $`A_n`$-modules. This was the 2004-2005 question for our algebraic combinatorics working seminar at Fields Institute where the research for this article was done.
Our answer involves the partition lattice algebras $`(\mathrm{𝕜}\mathrm{\Pi }_n,)`$ and $`(\mathrm{𝕜}\mathrm{\Pi }_n,)`$ (as well as the Solomon-Tits algebras ). For each one, with finite modules we can define a tensor product of $`\mathrm{𝕜}\mathrm{\Pi }_n`$ modules and a restriction from $`\mathrm{𝕜}\mathrm{\Pi }_n`$ module to $`\mathrm{𝕜}\mathrm{\Pi }_k\mathrm{𝕜}\mathrm{\Pi }_{nk}`$ modules. This allows us to place on $`_nG_0(\mathrm{𝕜}\mathrm{\Pi }_n)`$, the Grothendieck ring of the $`\mathrm{𝕜}\mathrm{\Pi }_n`$, a bialgebra structure (but not a Hopf algebra structure). We then define a bialgebra isomorphism $`_nG_0(\mathrm{𝕜}\mathrm{\Pi }_n)\mathrm{𝖭𝖢𝖲𝗒𝗆}^{}`$. We call this map the Frobenius characteristic map of the partition lattice algebras. This singles out a unique canonical basis of $`\mathrm{𝖭𝖢𝖲𝗒𝗆}`$ (up to automorphism) corresponding to the simple modules of the $`\mathrm{𝕜}\mathrm{\Pi }_n`$.
Our paper is divided into 4 sections as follows. In section 1 we recall the definition and structure of $`\mathrm{𝖭𝖢𝖲𝗒𝗆}`$. We then state our first theorem claiming the existence of a basis $`𝐱`$ of $`\mathrm{𝖭𝖢𝖲𝗒𝗆}`$ defined by certain algebraic properties. The proof of it will be postponed to section 4. In section 2 we recall the definition and structure of the partition lattice algebras $`\mathrm{𝕜}\mathrm{\Pi }_n`$ with the product given by the lattice operation $``$ and define their modules. We then introduce a structure of a semi-tower of algebras (i.e. we have a non-unital embedding $`\rho _{n,m}:\mathrm{𝕜}\mathrm{\Pi }_n\mathrm{𝕜}\mathrm{\Pi }_m\mathrm{𝕜}\mathrm{\Pi }_{n+m}`$ of algebras) on the partition lattice algebras and show that it induces a bialgebra structure on its Grothendieck ring. Our second theorem states that this Grothendieck bialgebra is dual to $`\mathrm{𝖭𝖢𝖲𝗒𝗆}`$. The classes of simple modules correspond then to the basis $`𝐱`$. In view of the work of Brown we remark that this can also be done with the semi-tower of Solomon-Tits algebras. In section 3 we build the same construction with the lattice algebras $`\mathrm{𝕜}\mathrm{\Pi }_n`$ with the product $``$. With this tower of algebras (i.e. $`\rho _{n,m}`$ is a unital morphism of algebras) we find that the Grothendieck bialgebra is again dual to $`\mathrm{𝖭𝖢𝖲𝗒𝗆}`$, but this time the classes of simple modules correspond to the monomial basis of $`\mathrm{𝖭𝖢𝖲𝗒𝗆}`$.
In section 4 we give the proof of our first theorem and show the basis canonically defined in section 2 corresponds to the simple modules of the $`\mathrm{𝕜}\mathrm{\Pi }_n`$. In light of the Frobenius characteristic of section 2, the basis can be interpreted as an analogue of the Schur functions for $`\mathrm{𝖭𝖢𝖲𝗒𝗆}`$ and providing an answer to an open question of .
## 1. $`\mathrm{𝖭𝖢𝖲𝗒𝗆}`$ and the basis $`\{𝐱_A\}`$
We recall the basic definition and structure of $`\mathrm{𝖭𝖢𝖲𝗒𝗆}`$. Most of it can be found in . A set partition $`A`$ of $`m`$ is a set of non-empty subsets $`A_1,A_2,\mathrm{},A_k[m]=\{1,2,\mathrm{},m\}`$ such that $`A_iA_j=\mathrm{}`$ for $`ij`$ and $`A_1A_2\mathrm{}A_k=[m]`$. The subsets $`A_i`$ are called the parts of the set partition and the number of non-empty parts the length of $`A`$, denoted by $`\mathrm{}(A)`$. There is a natural mapping from set partitions to integer partitions given by $`\lambda (A)=(|A_1|,|A_2|,\mathrm{},|A_k|)`$, where we assume that the blocks of the set partition have been listed in weakly decreasing order of size. We shall use $`\mathrm{}(\lambda )`$ to refer to the length (the number of parts) of the partition and $`|\lambda |`$ is the size of the partition (the sum of the sizes of the parts), while $`n_i(\lambda )`$ shall refer to the number of parts of the partition of size $`i`$. We denote by $`\mathrm{\Pi }_m`$ the set of set partitions of $`m`$. The number of set partitions is given by the Bell numbers. These can be defined by the recurrence $`B_0=1`$ and $`B_n=_{i=0}^{n1}\left(\genfrac{}{}{0pt}{}{n1}{i}\right)B_i`$.
For a set $`S=\{s_1,s_2,\mathrm{},s_k\}`$ of integers $`s_i`$ and an integer $`n`$ we use the notation $`S+n`$ to represent the set $`\{s_1+n,s_2+n,\mathrm{},s_k+n\}`$. For $`A\mathrm{\Pi }_m`$ and $`B\mathrm{\Pi }_r`$ set partitions with parts $`A_i`$, $`1i\mathrm{}(A)`$ and $`B_i`$, $`1i\mathrm{}(B)`$ respectively, we set $`A|B=\{A_1,A_2,\mathrm{},A_{\mathrm{}(A)},B_1+m,B_2+m,\mathrm{},B_{\mathrm{}(B)}+m\}`$, therefore $`A|B\mathrm{\Pi }_{m+r}`$ and this operation is noncommutative in the sense that, in general, $`A|BB|A`$.
When writing examples of set partitions, whenever the context allows it, we will use a more compact notation. For example, $`\{\{1,3,5\},\{2\},\{4\}\}`$ will be represented by $`\{\mathrm{135.2.4}\}`$. Although there is no order on the parts of a set partition, we will impose an implied order such that the parts are arranged by increasing value of the smallest element in the subset. This implied order will allow us to reference the $`i^{th}`$ parts of the set partition without ambiguity.
There is a natural lattice structure on the set partitions of a given $`n`$. We define for $`A,B\mathrm{\Pi }_n`$ that $`AB`$ if for each $`A_iA`$ there is a $`B_jB`$ such that $`A_iB_j`$ (otherwise stated, that $`A`$ is finer than $`B`$). The set of set partitions of $`[n]`$ with this order forms a poset with rank function given by $`n`$ minus the length of the set partition. This poset has a unique minimal element $`\mathrm{𝟎}_n=\{1.2.\mathrm{}.n\}`$ and a unique maximal element $`\mathrm{𝟏}_n=\{12\mathrm{}n\}`$. The largest element smaller than both $`A`$ and $`B`$ is denoted $`AB=\{A_iB_j:1i\mathrm{}(A),1j\mathrm{}(B)\}`$ while the smallest element larger than $`A`$ and $`B`$ is denoted $`AB`$. The lattice $`(\mathrm{\Pi }_n,,)`$ is called the partition lattice.
###### Example 1.1.
Let $`A=\{\mathrm{138.24.5.67}\}`$ and $`B=\{\mathrm{1.238.4567}\}`$. $`A`$ and $`B`$ are not comparable in the inclusion order on set partitions. We calculate that $`AB=\{\mathrm{1.2.38.4.5.67}\}`$ and $`AB=\{12345678\}`$.
When a collection of disjoint sets of positive integers is not a set partition because the union of the parts is not $`[n]`$ for some $`n`$, we may lower the values in the sets so that they keep their relative values so that the resulting collection is a set partition (of an $`m<n`$). This operation is referred to as the ‘standardization’ of a set of disjoint sets $`A`$ and the resulting set partition is denoted $`st(A)`$.
Now for $`A\mathrm{\Pi }_m`$ and $`S\{1,2,\mathrm{},\mathrm{}(A)\}`$ with $`S=\{s_1,s_2,\mathrm{},s_k\}`$, we define $`A_S=st(\{A_{s_1},A_{s_2},\mathrm{},A_{s_k}\})`$ which is a set partition of $`|A_{s_1}|+|A_{s_2}|+\mathrm{}+|A_{s_k}|`$. By convention $`A_{\{\}}`$ is the empty set partition.
###### Example 1.2.
If $`A=\{\mathrm{1368.2.4.579}\}`$, then $`A_{\{1,4\}}=\{1246.357\}`$.
For $`n0`$, consider a set $`X_n`$ of non-commuting variables $`x_1,x_2,\mathrm{},x_n`$ and the polynomial algebra $`_{X_n}=\mathrm{𝕜}x_1,x_2,\mathrm{},x_n`$ in these non-commuting variables. There is a natural $`S_n`$ action on the basis elements defined by $`\sigma (x_{i_1}x_{i_2}\mathrm{}x_{i_k})=x_{\sigma (i_1)}x_{\sigma (i_2)}\mathrm{}x_{\sigma (i_k)}.`$ Let $`x_{i_1}x_{i_2}\mathrm{}x_{i_m}`$ be a monomial in the space $`_{X_n}`$. We say that the type of this monomial is a set partition $`A\mathrm{\Pi }_m`$ with the property that $`i_a=i_b`$ if and only if $`a`$ and $`b`$ are in the same block of the set partition. This set partition is denoted as $`(i_1,i_2,\mathrm{},i_m)=A`$. Notice that the length of $`(i_1,i_2,\mathrm{},i_m)`$ is equal to the number of different values which appear in $`(i_1,i_2,\mathrm{},i_m)`$.
The vector space $`\mathrm{𝖭𝖢𝖲𝗒𝗆}^{(n)}`$ is defined as the linear span of the elements
$$\text{m}_A[X_n]=\underset{(i_1,i_2,\mathrm{},i_m)=A}{}x_{i_1}x_{i_2}\mathrm{}x_{i_m}$$
for $`A\mathrm{\Pi }_m`$, where the sum is over all sequences with $`1i_jn`$. For the empty set partition, we define by convention $`\text{m}_{\{\}}[X_n]=1`$. If $`\mathrm{}(A)>n`$ we must have that $`\text{m}_A[X_n]=0`$. Since for any permutation $`\sigma S_m`$, $`(i_1,i_2,\mathrm{},i_m)=(\sigma (i_1),\sigma (i_2),\mathrm{},\sigma (i_m))`$, we have that $`\sigma \text{m}_A[X_n]=\text{m}_A[X_n]`$. In fact, $`\text{m}_A[X_n]`$ is the sum of all elements in the orbit of a monomial of type $`A`$ under the action of $`S_n`$. Therefore $`\mathrm{𝖭𝖢𝖲𝗒𝗆}^{(n)}`$ is the space of $`S_n`$-invariants in the noncommutative polynomial algebra $`_{X_n}`$. For instance, $`\text{m}_{\{13.2\}}[X_4]=x_1x_2x_1+x_1x_3x_1+x_1x_4x_1+x_2x_1x_2+x_2x_3x_2+x_2x_4x_2+x_3x_1x_3+x_3x_2x_3+x_3x_4x_3+x_4x_1x_4+x_4x_2x_4+x_4x_3x_4.`$
Now let $`\mathrm{𝖭𝖢𝖲𝗒𝗆}`$ be the inverse limit of $`\mathrm{𝖭𝖢𝖲𝗒𝗆}^{(n)}`$. We have that $`\mathrm{𝖭𝖢𝖲𝗒𝗆}=_{d0}\mathrm{𝖭𝖢𝖲𝗒𝗆}_d`$ is a graded algebra where $`\mathrm{𝖭𝖢𝖲𝗒𝗆}_d`$ is the linear span of $`\{\text{m}_A\}_{A\mathrm{\Pi }_d}`$. Here we forget any reference to the variables $`x_1,x_2,\mathrm{}`$ and think of elements in $`\mathrm{𝖭𝖢𝖲𝗒𝗆}`$ as noncommutative symmetric functions. The degree of a basis element $`\text{m}_A`$ is given by $`|A|=d`$ and the product map $`\mu :\mathrm{𝖭𝖢𝖲𝗒𝗆}_d\mathrm{𝖭𝖢𝖲𝗒𝗆}_m\mathrm{𝖭𝖢𝖲𝗒𝗆}_{d+m}`$ is defined on the basis elements $`\text{m}_A\text{m}_B`$ by
(1.1)
$$\mu (\text{m}_A\text{m}_B):=\underset{\genfrac{}{}{0pt}{}{C\mathrm{\Pi }_{d+m}}{C\mathbf{\hspace{0.17em}\hspace{0.17em}1}_d|\mathrm{𝟏}_m=A|B}}{}\text{m}_C$$
This is a lift of the multiplication in $`\mathrm{𝖭𝖢𝖲𝗒𝗆}^{(n)}`$ as $`n\mathrm{}`$.
The graded algebra $`\mathrm{𝖭𝖢𝖲𝗒𝗆}`$ is in fact a Hopf algebra with the following comultiplication $`\mathrm{\Delta }:\mathrm{𝖭𝖢𝖲𝗒𝗆}_d_{k=0}^d\mathrm{𝖭𝖢𝖲𝗒𝗆}_k\mathrm{𝖭𝖢𝖲𝗒𝗆}_{dk}`$ where
(1.2)
$$\mathrm{\Delta }(\text{m}_A)=\underset{S[\mathrm{}(A)]}{}\text{m}_{A_S}\text{m}_{A_{S^c}}$$
and $`S^c=[\mathrm{}(A)]S`$. The counit is given by $`ϵ:\mathrm{𝖭𝖢𝖲𝗒𝗆}`$ where $`ϵ(\text{m}_{\{\}})=1`$ and $`ϵ(\text{m}_A)=0`$ for all $`A\mathrm{\Pi }_n`$ for $`n>0`$. More details on this Hopf algebra structure are found in . The algebra $`\mathrm{𝖭𝖢𝖲𝗒𝗆}`$ was originally considered by Wolf in extending the fundamental theorem of symmetric functions to this algebra and later by Bergman and Cohn . More recently Rosas and Sagan considered this space to define natural bases which are analogous to bases of the (commutative) symmetric functions. More progress in understanding this space was made in where it was considered as a Hopf algebra. In the Hopf algebra $`Sym`$ of (commutative) symmetric functions, the comultiplication corresponds to the plethysm $`f[X]f[X+Y]`$. It was established in that the comultiplication in $`\mathrm{𝖭𝖢𝖲𝗒𝗆}`$ corresponds to a noncommutative plethysm $`F[X]F[X+Y]`$.
The Hopf algebra $`Sym`$ has more structure. There is a second comultiplication corresponding to the plethysm $`f[X]f[XY]`$ (see ). This second operation is often refer to as the internal comultiplication or Kronecker comultiplication. We end this section describing for $`\mathrm{𝖭𝖢𝖲𝗒𝗆}`$ the analog of this internal comultiplication. This description is also considered in .
For the Hopf algebra $`\mathrm{𝖭𝖢𝖲𝗒𝗆}`$ we define a second (internal) comultiplication $`\mathrm{\Delta }^{}:\mathrm{𝖭𝖢𝖲𝗒𝗆}_d\mathrm{𝖭𝖢𝖲𝗒𝗆}_d\mathrm{𝖭𝖢𝖲𝗒𝗆}_d`$ given by
(1.3)
$$\mathrm{\Delta }^{}(\text{m}_A)=\underset{BC=A}{}\text{m}_B\text{m}_C$$
This operation corresponds to a noncommutative plethysm $`F[X]F[XY]`$. More precisely, we first assume that we have two countable alphabet (totally ordered noncommutative variables) $`X=x_1,x_2,\mathrm{}`$ and $`Y=y_1,y_2,\mathrm{}`$, and we let $`XY=\{x_iy_j\}`$ be the alphabet totally ordered by the lexicographic order on the couple $`(i,j)`$. The transformation $`F[X]F[XY]`$ denotes the substitution given by the (countable) bijection $`XXY`$. If we let the $`x_i`$’s commute with the $`y_j`$’s then we have that $`F[XY]`$ can be expanded in the form $`F[XY]=F_{1,i}[X]F_{2,i}[Y]`$. We can then define the operation $`\mathrm{\Delta }^{}(F)=F_{1,i}F_{2,i}`$. In equation (1.3) we have given the result of this when $`F=\text{m}_A`$. Clearly this operation is a morphism for the multiplication, thus $`\mathrm{𝖭𝖢𝖲𝗒𝗆}`$ with $`\mathrm{\Delta }^{}`$ and the multiplication operation of equation (1.1) forms a bialgebra. But it is not a Hopf algebra as it does not have an antipode (the equation $`\mu (\text{m}_{\{\}}S)\mathrm{\Delta }^{}(\text{m}_{\{1\}})=\text{m}_{\{\}}`$ has no solutions, where $`S`$ would have been the antipode). We are now in position to state our first main theorem.
###### Theorem 1.3.
There is a basis $`\{𝐱_A:A\mathrm{\Pi }_n,n0\}`$ of $`\mathrm{𝖭𝖢𝖲𝗒𝗆}`$ such that
$`\text{(i)}𝐱_A𝐱_B=𝐱_{A|B}`$
$`\text{(ii)}\mathrm{\Delta }^{}(𝐱_C)={\displaystyle \underset{AB=C}{}}𝐱_A𝐱_B.`$
The proof of this theorem is technical and we differ it to Section 4. We are convinced that the basis $`\{𝐱_A:A\mathrm{\Pi }_n,n0\}`$ is central in the study of $`\mathrm{𝖭𝖢𝖲𝗒𝗆}`$ and should have many fascinating properties. We plan to study this basis further in future work. For now, we prefer to develop the representation theory that will motivate our result.
## 2. Grothendieck bialgebras of the Semi-tower $`(𝚷,)=_{n0}(\mathrm{𝕜}\mathrm{\Pi }_n,)`$.
In this section we consider the partition lattice algebras. For a fixed $`n`$ consider the vector space $`(\mathrm{𝕜}\mathrm{\Pi }_n,)`$ formally spanned by the set partitions of $`n`$. The multiplication is given by the operation $``$ on set partitions and with the unit $`\mathrm{𝟏}_n=\{[n]\}`$. We remark that for all $`d`$, we have that $`\mathrm{𝕜}\mathrm{\Pi }_d`$ is isomorphic as a vector space to $`\mathrm{𝖭𝖢𝖲𝗒𝗆}_d`$ via the pairing $`A\text{m}_A`$. Moreover, it is straightforward to check using equation (1.3) that $`\mathrm{\Delta }^{}`$ is dual to $``$ as operators.
It is well known that $`(\mathrm{𝕜}\mathrm{\Pi }_n,)`$ is a commutative semisimple algebra (see \[19, Theorem 3.9.2\]). To see this, one considers the algebra $`\mathrm{𝕜}^{\mathrm{\Pi }_n}=\{f:\mathrm{\Pi }_n\mathrm{𝕜}\}`$ which is clearly commutative and semisimple. We then define the map
$`\delta _{}:(\mathrm{𝕜}\mathrm{\Pi }_n,)`$ $`\mathrm{𝕜}^{\mathrm{\Pi }_n}`$
$`A`$ $`\delta _A,`$
where $`\delta _A(B)=1`$ if $`AB`$ and $`0`$ otherwise. Next check that $`\delta _{AB}=\delta _A\delta _B`$ which shows that $`\delta _{}`$ is an isomorphism of algebras.
The primitive orthogonal idempotents of $`\mathrm{𝕜}^{\mathrm{\Pi }_n}`$ are given by the functions $`\delta _{A=}`$ defined by $`\delta _{A=}(B)=1`$ if $`A=B`$ and $`0`$ otherwise. We have that $`\delta _A=_{BA}\delta _{=B}`$. This implies, using Möbius inversion, that the primitive orthogonal idempotents of $`(\mathrm{𝕜}\mathrm{\Pi }_n,)`$ are given by
(2.1)
$$e_A=\underset{BA}{}\mu (B,A)B,$$
where $`\mu `$ is the Möbius function of the partially ordered set $`\mathrm{\Pi }_n`$. Since $`(\mathrm{𝕜}\mathrm{\Pi }_n,)`$ is commutative and semisimple, we have that the simple $`(\mathrm{𝕜}\mathrm{\Pi }_n,)`$-modules of this algebra are the one dimensional spaces $`V_A=\mathrm{𝕜}\mathrm{\Pi }_ne_A`$. Here the action is given by the left multiplication
(2.2)
$$Ce_A=\{\genfrac{}{}{0pt}{}{e_A\text{if }CA\text{,}}{0\text{otherwise.}}$$
This follows from the corresponding identity in $`\mathrm{𝕜}^{\mathrm{\Pi }_n}`$ considering $`\delta _C\delta _{=A}`$.
We now let $`G_0(\mathrm{𝕜}\mathrm{\Pi }_n,)`$ denote the Grothendieck group of the category of finite dimensional $`(\mathrm{𝕜}\mathrm{\Pi }_n,)`$-modules. This is the vector space spanned by the equivalence classes of simple $`(\mathrm{𝕜}\mathrm{\Pi }_n,)`$-modules under isomorphisms. We also consider $`K_0(\mathrm{𝕜}\mathrm{\Pi }_n,)`$ the Grothendieck group of the category of projective $`(\mathrm{𝕜}\mathrm{\Pi }_n,)`$-modules. Since $`(\mathrm{𝕜}\mathrm{\Pi }_n,)`$ is semisimple, the space $`G_0(\mathrm{𝕜}\mathrm{\Pi }_n,)`$ and $`K_0(\mathrm{𝕜}\mathrm{\Pi }_n,)`$ are equal as vector spaces as they are both linearly spanned by the elements $`V_A`$ for $`A\mathrm{\Pi }_n`$. We then set $`K_0(𝚷,)=_{n0}K_0(\mathrm{𝕜}\mathrm{\Pi }_n,)`$.
Given two finite $`(\mathrm{𝕜}\mathrm{\Pi }_n,)`$ modules $`V`$ and $`W`$, we can form the $`(\mathrm{𝕜}\mathrm{\Pi }_n,)`$-module $`VW`$ with the diagonal action (it is an action since a semigroup algebra is a bialgebra for the coproduct $`AAA`$). We denote this $`(\mathrm{𝕜}\mathrm{\Pi }_n,)`$-module by $`VW`$ (to avoid confusion with the tensor product of a $`(\mathrm{𝕜}\mathrm{\Pi }_n,)`$-module and a $`(\mathrm{𝕜}\mathrm{\Pi }_m,)`$-module).
###### Lemma 2.1.
Given two simple $`(\mathrm{𝕜}\mathrm{\Pi }_n,)`$-module $`V_A`$ and $`V_B`$,
(2.3)
$$V_AV_B=V_{AB}.$$
###### Proof.
Let $`C\mathrm{\Pi }_n`$ act on $`e_Ae_B`$. ¿From equation (2.2) we get $`C(e_Ae_B)=(Ce_A)(Ce_B)=e_Ae_B`$ if and only if $`CA`$ and $`CB`$, that is $`CAB`$. If not, we get $`C(e_Ae_B)=0`$. We conclude that the map $`e_Ae_Be_{AB}`$ is the desired isomorphism in equation (2.3). ∎
We would like to define on $`G_0(𝚷,)=_{n0}G_0(\mathrm{𝕜}\mathrm{\Pi }_n,)`$ a graded multiplication and a graded comultiplication corresponding to induction and restriction. For this we need a few more tools.
###### Lemma 2.2.
The linear map $`\rho _{n,m}:(\mathrm{𝕜}\mathrm{\Pi }_n,)(\mathrm{𝕜}\mathrm{\Pi }_m,)(\mathrm{𝕜}\mathrm{\Pi }_{n+m},)`$ defined by $`\rho _{n,m}(AB)=A|B`$ is injective and multiplicative. Moreover, $`\rho _{k+n,m}(\rho _{k,n}Id)=\rho _{k,n+m}(Id\rho _{n,m})`$ for all $`k,n`$ and $`m`$.
###### Proof.
Let $`A=\{A_1,\mathrm{},A_r\}`$, $`B=\{B_1,\mathrm{},B_s\}`$, $`C=\{C_1,\mathrm{},C_t\}`$ and $`D=\{D_1,\mathrm{},D_u\}`$, where $`A,B\mathrm{\Pi }_n`$ and $`C,D\mathrm{\Pi }_m`$. We remark that for all $`i,j`$, we have $`A_i(D_j+n)=\mathrm{}`$ and $`(C_i+n)B_j=\mathrm{}`$. Since $`(C_i+n)(D_j+n)=(C_iD_j)+n`$, we have
$`(A|C)(B|D)`$ $`=\left\{A_iB_j\right\}_{\genfrac{}{}{0pt}{}{1ir}{1js}}\left\{(C_i+n)(D_j+n)\right\}_{\genfrac{}{}{0pt}{}{1it}{1ju}}`$
$`=(AB)|(CD),`$
and this shows that $`\rho _{n,m}`$ is multiplicative. The injectivity of this map is clear from the fact that $`\rho _{n,m}`$ maps distinct basis elements into distinct basis elements. The last identity of the lemma follows from the associativity of the operation “$`|`$” ∎
We define a semi-tower $`(_{n0}A_n,\{\varphi _{n,m}\})`$ to be a direct sum of algebras along with a family of injective non-unital homomorphisms of algebras $`\varphi _{n,m}:A_nA_mA_{n+m}`$. A tower in the sense defined in the recent literature is a semi-tower with the additional constraint that $`\varphi _{n,m}(\mathrm{𝟏}_n,\mathrm{𝟏}_m)=\mathrm{𝟏}_{n+m}`$ (i.e. that $`\varphi _{n,m}`$ is a unital embedding of algebras).
Define the pair $`(𝚷,)=(_{n0}(\mathrm{𝕜}\mathrm{\Pi }_n,),\{\rho _{n,m}\})`$ which is a semi-tower of the algebras $`(\mathrm{𝕜}\mathrm{\Pi }_n,)`$. We remark that $`(𝚷,)`$ is a graded algebra with the multiplication $`\rho _{n,m}(A,B)=A|B`$ which is associative (but non-commutative) and has a unit given by the emptyset partition $`\mathrm{}\mathrm{\Pi }_0`$. Moreover, each of the homogeneous components $`(\mathrm{𝕜}\mathrm{\Pi }_n,)`$ of $`𝚷`$ are themselves algebras with the multiplication $``$, and Lemma 2.2 gives the relationship between the two operations.
At this point we need to stress that $`\rho _{n,m}`$ is not a unital embedding of algebras and hence $`(𝚷,)`$ is not a tower of algebras. The algebra $`(\mathrm{𝕜}\mathrm{\Pi }_n,)`$ has a unit given by $`\mathrm{𝟏}_n=\{12\mathrm{}n\}`$, but $`\rho _{n,m}(\mathrm{𝟏}_n\mathrm{𝟏}_m)\mathrm{𝟏}_{n+m}`$. The tower of algebras considered in the recent literature all have the property that the corresponding $`\rho _{n,m}`$ are (unital) embeddings of algebras. This is the reason we call our construction a semi-tower rather than a tower.
The motivation for defining a tower of algebras is to allow one to induce and restrict modules of these algebras and ultimately to define on its Grothendieck ring a Hopf algebra structure. Here the fact that we have only a semi-tower causes some problems in defining restriction of modules. Yet we can still define a weaker version of restriction in our situation. Let $`A`$ and $`B`$ be two finite dimensional algebras and let $`\rho :AB`$ be a multiplicative injective linear map. Given a finite $`B`$-module $`M`$, we define
$$\text{Res}_\rho M=\{mM:\rho (\mathrm{𝟏}_A)m=m\}M.$$
In the case where $`\rho `$ is an embedding of algebras this definition agrees with the traditional one. More on this general theory will be found in but here we focus our attention on $`(𝚷,)`$.
###### Lemma 2.3.
For $`kn`$ and a simple $`(\mathrm{𝕜}\mathrm{\Pi }_n,)`$-module $`V_AG_0(\mathrm{𝕜}\mathrm{\Pi }_n,)`$,
$$\text{Res}_{\rho _{k,nk}}V_A=\{\begin{array}{cc}V_A\hfill & \text{if }A=B|C\text{ for }B\mathrm{\Pi }_k\text{ and }C\mathrm{\Pi }_{nk}\text{ }\hfill \\ 0\hfill & \text{ otherwise.}\hfill \end{array}$$
###### Proof.
We have that $`\rho _{n,m}(\mathrm{𝟏}_k\mathrm{𝟏}_{nk})e_A=(\mathrm{𝟏}_k|\mathrm{𝟏}_{nk})e_A=e_A`$ if $`\mathrm{𝟏}_k|\mathrm{𝟏}_{nk}A`$, and $`0`$ otherwise. The condition $`\mathrm{𝟏}_k|\mathrm{𝟏}_{nk}A`$ is equivalent to $`A=B|C`$ where $`A|_{1,\mathrm{},k}=B`$ and $`A|_{k+1,\mathrm{},n+k}=C`$. ∎
We can now define a graded comultiplication on $`G_0(𝚷,)`$ using our definition of restriction. For $`VG_0(\mathrm{𝕜}\mathrm{\Pi }_n,)`$ let
(2.4)
$$\mathrm{\Delta }(V)=\underset{k=0}{\overset{n}{}}\text{Res}_{\rho _{k,nk}}V.$$
It follows from Lemmas 2.2 that this operation is coassociative. For a simple module $`V_AG_0(\mathrm{𝕜}\mathrm{\Pi }_n,)`$, Lemma 2.3 gives us
(2.5)
$$\mathrm{\Delta }(V_A)=\underset{A=B|C}{}V_BV_C.$$
Now we extend $``$ to $`G_0(𝚷,)`$ by setting $`V_AV_B=0`$ if $`V_A`$ and $`V_B`$ are not of the same degree.
###### Proposition 2.4.
$`(G_0(𝚷,),,\mathrm{\Delta })`$ is a bialgebra.
###### Proof.
Let $`A,B\mathrm{\Pi }_n`$. By equation (2.3), it is sufficient to prove that $`\mathrm{\Delta }(V_{AB})=\mathrm{\Delta }(V_A)\mathrm{\Delta }(V_B)`$. Using equation (2.3) we can easily reduce the problem to the following assertion: there are $`C\mathrm{\Pi }_k`$, $`D\mathrm{\Pi }_{nk}`$ such that $`AB=C|D`$ if and only if there are $`E,E^{}\mathrm{\Pi }_k`$, $`F,F^{}\mathrm{\Pi }_{nk}`$ such that $`A=E|F`$ and $`B=E^{}|F^{}`$. This follows then from definitions. ∎
It is thus natural to give a notion to induced modules dual to restriction in Lemma 2.3.
###### Lemma 2.5.
For two simple modules $`V_A=\mathrm{𝕜}\mathrm{\Pi }_ne_AG_0(\mathrm{𝕜}\mathrm{\Pi }_n,)`$ and $`V_B=\mathrm{𝕜}\mathrm{\Pi }_me_BG_0(\mathrm{𝕜}\mathrm{\Pi }_m,)`$ we define
$$\text{Ind}_{n,m}V_AV_B=\mathrm{𝕜}\mathrm{\Pi }_{n+m}_{\mathrm{𝕜}\mathrm{\Pi }_n\mathrm{𝕜}\mathrm{\Pi }_m}(\mathrm{𝕜}\mathrm{\Pi }_ne_A\mathrm{𝕜}\mathrm{\Pi }_me_B),$$
where $`\mathrm{𝕜}\mathrm{\Pi }_n\mathrm{𝕜}\mathrm{\Pi }_m`$ is embedded into $`\mathrm{𝕜}\mathrm{\Pi }_{n+m}`$ via $`\rho _{n,m}`$.
There is a natural isomorphism such that
$$\text{Ind}_{n,m}V_AV_B\mathrm{𝕜}\mathrm{\Pi }_{n+m}\rho _{n,m}(e_Ae_B).$$
We have
(2.6)
$$\text{Ind}_{n,m}V_AV_B=V_{A|B}.$$
###### Proof.
Consider the following isomorphism which allows us to naturally realize $`\text{Ind}_{n,m}V_AV_B`$ as an element of $`G_0(\mathrm{𝕜}\mathrm{\Pi }_{n+m},)`$.
$`\text{Ind}_{n,m}V_AV_B`$ $`=\mathrm{𝕜}\mathrm{\Pi }_{n+m}_{\mathrm{𝕜}\mathrm{\Pi }_n\mathrm{𝕜}\mathrm{\Pi }_m}(\mathrm{𝕜}\mathrm{\Pi }_ne_A\mathrm{𝕜}\mathrm{\Pi }_me_B)`$
$`=\mathrm{𝕜}\mathrm{\Pi }_{n+m}_{\mathrm{𝕜}\mathrm{\Pi }_n\mathrm{𝕜}\mathrm{\Pi }_m}(e_Ae_B)`$
$`=\mathrm{𝕜}\mathrm{\Pi }_{n+m}\rho _{n,m}(e_Ae_B)_{\mathrm{𝕜}\mathrm{\Pi }_n\mathrm{𝕜}\mathrm{\Pi }_m}(\mathrm{𝟏}_n\mathrm{𝟏}_n)`$
$`\mathrm{𝕜}\mathrm{\Pi }_{n+m}\rho _{n,m}(e_Ae_B).`$
By linearity
$$\rho _{n,m}(e_Ae_B)=e_A|e_B=\underset{CA}{}\underset{DB}{}\mu (C,A)\mu (D,B)C|D.$$
We now remark that $`\{E:EA|B\}=\{C|D:C|DA|B\}=\{C|D:CA,DB\}`$. This is isomorphic to the cartesian product $`\{C:CA\}\times \{D:DB\}`$. Since Möbius functions are multiplicative with respect to cartesian product we have
$$\rho _{n,m}(e_Ae_B)=\underset{EA|B}{}\mu (E,A|B)E=e_{A|B}.$$
It is clear now that $`\text{Ind}_{n,m}`$ defines on $`G_0(𝚷,)`$ a graded multiplication $`V_AV_BV_{A|B}`$ that is dual to the graded comultiplication of $`\mathrm{\Delta }`$ defined on $`G_0(𝚷,)`$. We also define an internal comultiplication on $`G_0(𝚷,)`$ dual to equation (2.3) such that $`\mathrm{\Delta }^{}:G_0(\mathrm{𝕜}\mathrm{\Pi }_n,)G_0(\mathrm{𝕜}\mathrm{\Pi }_n,)G_0(\mathrm{𝕜}\mathrm{\Pi }_n,)`$. For $`C\mathrm{\Pi }_n`$ let
(2.7)
$$\mathrm{\Delta }^{}(V_C)=\underset{AB=C}{}V_AV_B.$$
The space $`G_0(𝚷,)`$ with its graded multiplication given by induction and comultiplication $`\mathrm{\Delta }^{}`$ is a bialgebra, by duality and Proposition 2.4. The main theorem of this section is a direct corollary to Theorem 1.3.
###### Theorem 2.6.
The map $`F:G_0(𝚷,)\mathrm{𝖭𝖢𝖲𝗒𝗆}`$ defined by $`F(V_A)=𝐱_A`$ is an isomorphism of bialgebras.
###### Proof.
$`G_0(𝚷,)`$ is endowed with a product given by (2.6) and an inner coproduct given by (2.7). Since $`\mathrm{𝖭𝖢𝖲𝗒𝗆}`$ is known to be a bialgebra satisfying the relations given in Theorem 1.3, the map $`F`$ is an isomorphism. ∎
The map $`F`$ is called the Frobenius map for our semi-tower. Along with Theorem 1.3, it shows that the basis $`𝐱_A`$ of $`\mathrm{𝖭𝖢𝖲𝗒𝗆}`$ are the only functions that correspond to the classes of simple modules in $`G_0(𝚷,)`$. This defines $`𝐱_A`$ uniquely (up to automorphism) and for this reason we think of them as the Schur functions for the semi-tower $`(𝚷,)`$ of the symmetric functions in non-commutative variables.
###### Remark 2.7.
In , Brown shows that $`(\mathrm{𝕜}\mathrm{\Pi }_n,)`$ is the semisimple quotient of the Solomon-Tits algebra $`ST_n`$ (see ). It is easy to lift our semi-tower structure from $`𝚷`$ to $`\mathrm{𝐒𝐓}=ST_n`$ via Brown’s support map. Then $`G_0(\mathrm{𝐒𝐓},)`$ and $`G_0(𝚷,)`$ are isomorphic as bialgebras.
In , the conditions under which a tower of algebras $`𝐀=(_{n0}A_n,\rho _{n,m})`$ defines a Hopf algebra structure on the Grothendieck rings $`G_0(𝐀)`$ and $`K_0(𝐀)`$ are considered. Under certain conditions one would expect that the Grothendieck ring $`G_0(𝐀)`$ of finite modules forms a Hopf algebra with the operations of induction and restriction which is isomorphic to the graded dual of the Grothendieck ring $`K_0(𝐀)`$ of projective modules.
For the tower of algebras we are considering here, it is not the case that $`G_0(𝚷,)`$ forms a Hopf algebra because the operations of induction and restriction are not even compatible as a bialgebra structure. We have shown that $`G_0(𝚷,)`$ and $`K_0(𝚷,)`$ are endowed naturally with a product given by the notion of induction in equation (2.6) and coproduct given by the notion of restriction given in equation (2.4). It is easily checked that these operations do not form a Hopf algebra structure.
We have found however that here we have $`G_0(𝚷,)`$ endowed with the operations of induction and restriction is isomorphic by the graded dual to $`K_0(𝚷,)`$ also endowed with the same induction and restriction operations. This is because the operation of restriction on $`G_0(𝚷,)`$ is dual to the operation of induction on $`K_0(𝚷,)`$ and induction on $`G_0(𝚷,)`$ is dual as graded operations to restriction on $`K_0(𝚷,)`$. This remark can be observed through the duality in equations (2.6) and (2.4).
## 3. Grothendieck bialgebras of the Tower $`(𝚷,)=_{n0}(\mathrm{𝕜}\mathrm{\Pi }_n,)`$.
In this section we consider a second algebra related to the partition lattice and show that there is an additional connection with the algebra $`\mathrm{𝖭𝖢𝖲𝗒𝗆}`$. Define $`(\mathrm{𝕜}\mathrm{\Pi }_n,)`$ to be the commutative algebra linearly spanned by the elements of $`\mathrm{\Pi }_n`$ and endowed with the product $``$. This algebra has as a unit the minimal element $`\mathrm{𝟎}_n=\{1.2.\mathrm{}.n\}`$ of the poset $`\mathrm{\Pi }_n`$ since $`\mathrm{𝟎}_nA=A`$ for all $`A\mathrm{\Pi }_n`$.
As we constructed the primitive orthogonal idempotents for $`(\mathrm{𝕜}\mathrm{\Pi }_n,)`$, we proceed by defining in a similar manner
$`\delta _{}:(\mathrm{𝕜}\mathrm{\Pi }_n,)`$ $`\mathrm{𝕜}^{\mathrm{\Pi }_n}`$
$`A`$ $`\delta _A.`$
It is straightforward to check that $`\delta _A\delta _B=\delta _{(AB)}`$ and hence $`\delta _{}`$ is an isomorphism of algebras. This map can be used to recover the primitive orthogonal idempotents of $`(\mathrm{𝕜}\mathrm{\Pi }_n,)`$ since if $`\delta _A=_{BA}\delta _{B=}`$, then $`\delta _{A=}=_{BA}\mu (A,B)\delta _B`$. This can be summarized in the following proposition.
###### Proposition 3.1.
The primitive orthogonal idempotents of the algebra $`(\mathrm{𝕜}\mathrm{\Pi }_n,)`$ are
$$f_A=\underset{BA}{}\mu (A,B)B$$
with the property that
(3.1)
$$Cf_A=\{\begin{array}{cc}f_A\hfill & \text{if }CA\hfill \\ 0\hfill & \text{otherwise}\hfill \end{array}$$
It is also not difficult to check that the map $`\rho _{n,m}(A,B)=A|B`$ is also multiplicative with respect to the $``$ product in analogy with Lemma 2.2. Therefore we define the tower of algebras $`(𝚷,)=(_{n0}(\mathrm{𝕜}\mathrm{\Pi }_n,),\{\rho _{n,m}\})`$. This time we find that $`\rho _{n,m}`$ is indeed an embedding of algebras and $`(𝚷,)`$ a tower of algebras (see remarks related to $`(𝚷,)`$) since $`\rho _{n,m}(\mathrm{𝟎}_n,\mathrm{𝟎}_m)=\mathrm{𝟎}_{n+m}`$.
We now define $`G_0(\mathrm{𝕜}\mathrm{\Pi }_n,)`$ to be the ring of the category of finite dimensional $`(\mathrm{𝕜}\mathrm{\Pi }_n,)`$-modules endowed with the tensor of modules as the product. $`G_0(\mathrm{𝕜}\mathrm{\Pi }_n,)`$ is linearly spanned by the equivalence classes of the simple modules under isomorphism. Also set $`K_0(\mathrm{𝕜}\mathrm{\Pi }_n,)`$ to be the Grothendieck ring of the category of projective $`(\mathrm{𝕜}\mathrm{\Pi }_n,)`$-modules. Since $`(\mathrm{𝕜}\mathrm{\Pi }_n,)`$ is semi-simple we find that $`G_0(\mathrm{𝕜}\mathrm{\Pi }_n,)K_0(\mathrm{𝕜}\mathrm{\Pi }_n,)`$ and both are spanned by the simple modules $`W_A=\mathrm{𝕜}\mathrm{\Pi }_nf_A`$. Set $`G_0(𝚷,)=_{n0}G_0(\mathrm{𝕜}\mathrm{\Pi }_n,)`$ and $`K_0(𝚷,)=_{n0}K_0(\mathrm{𝕜}\mathrm{\Pi }_n,)`$.
Given two simple modules $`W_A,W_BG_0(\mathrm{𝕜}\mathrm{\Pi }_n,)`$ we consider the tensor product of modules $`W_AW_B`$ with the diagonal action of $`(\mathrm{𝕜}\mathrm{\Pi }_n,)`$. Denote this module as $`W_AW_B`$. We find that $`C(f_Af_B)=(Cf_A)(Cf_B)`$ which is equal to $`f_Af_B`$ if $`CA`$ and $`CB`$ (i.e. $`CAB`$) and it is equal to $`0`$ otherwise. We conclude from this discussion the following lemma.
###### Lemma 3.2.
For $`W_A,W_BG_0(\mathrm{𝕜}\mathrm{\Pi }_n,)`$,
(3.2)
$$W_AW_B=W_{AB}$$
is a simple module.
We have the following formula for the restriction of $`W_A`$ to $`\mathrm{𝕜}\mathrm{\Pi }_k\mathrm{𝕜}\mathrm{\Pi }_{nk}`$.
###### Lemma 3.3.
For $`kn`$ and a simple module $`W_AG_0(\mathrm{𝕜}\mathrm{\Pi }_n,)`$,
$$\text{Res}_{\rho _{k,nk}}W_A=W_BW_C$$
where $`A(\mathrm{𝟏}_k|\mathrm{𝟏}_{nk})=B|C`$ for $`B\mathrm{\Pi }_k`$ and $`C\mathrm{\Pi }_{nk}`$.
###### Proof.
First we check that $`\rho _{n,m}(\mathrm{𝟎}_k\mathrm{𝟎}_{nk})f_A=\mathrm{𝟎}_nf_A=f_A`$. Now for $`B^{}\mathrm{\Pi }_k`$ and $`C^{}\mathrm{\Pi }_{nk}`$, we have that $`\rho _{k,nk}(B^{},C^{})f_A=(B^{}|C^{})f_A=f_A`$ if $`(B^{}|C^{})A`$ and $`0`$ otherwise. If $`A(\mathrm{𝟏}_k|\mathrm{𝟏}_{nk})=(B|C)`$ then $`(B^{}|C^{})f_A=f_A`$ if and only if $`B^{}B`$ and $`C^{}C`$. Therefore $`W_A`$ is isomorphic to $`W_BW_C`$ as a $`\mathrm{𝕜}\mathrm{\Pi }_k\mathrm{𝕜}\mathrm{\Pi }_{nk}`$ module. ∎
Define now a notion of induction for $`K_0(𝚷,)`$ (as the dual of $`G_0(𝚷,)`$). For $`A\mathrm{\Pi }_n`$ and $`B\mathrm{\Pi }_m`$, the induced $`(\mathrm{𝕜}\mathrm{\Pi }_{n+m},)`$ module is
$$\text{Ind}_{n,m}W_AW_B=\mathrm{𝕜}\mathrm{\Pi }_{n+m}_{\mathrm{𝕜}\mathrm{\Pi }_n\mathrm{𝕜}\mathrm{\Pi }_m}(W_AW_B)$$
where we consider $`\mathrm{𝕜}\mathrm{\Pi }_n\mathrm{𝕜}\mathrm{\Pi }_m\rho _{n,m}(\mathrm{𝕜}\mathrm{\Pi }_n\mathrm{𝕜}\mathrm{\Pi }_m)\mathrm{𝕜}\mathrm{\Pi }_{n+m}`$.
###### Lemma 3.4.
For $`A\mathrm{\Pi }_n`$ and $`B\mathrm{\Pi }_m`$ we have that
(3.3)
$$\text{Ind}_{n,m}W_AW_B=\underset{C(\mathrm{𝟏}_n|\mathrm{𝟏}_m)=A|B}{}W_C.$$
###### Proof.
By proceeding as in the proof of Lemma 2.5, we get
$$\text{Ind}_{n,m}W_AW_B\mathrm{𝕜}\mathrm{\Pi }_{n+m}\rho _{n,m}(f_Af_B).$$
Therefore we just have to prove
$$\rho _{n,m}(f_Af_B)=\underset{C(\mathrm{𝟏}_n|\mathrm{𝟏}_m)=A|B}{}f_C.$$
Since Möbius functions are multiplicative with respect to cartesian product we have on the one hand
$$\rho _{n,m}(f_Af_B)=\underset{E|FA|B}{}\mu (A|B,E|F)E|F.$$
On the other hand
$`{\displaystyle \underset{C(\mathrm{𝟏}_n|\mathrm{𝟏}_m)=A|B}{}}f_C`$ $`=`$ $`{\displaystyle \underset{C(\mathrm{𝟏}_n|\mathrm{𝟏}_m)=A|B}{}}{\displaystyle \underset{DC}{}}\mu (C,D)D`$
$`=`$ $`{\displaystyle \underset{E|FA|B}{}}{\displaystyle \underset{C(\mathrm{𝟏}_n|\mathrm{𝟏}_m)=A|B}{}}{\displaystyle \underset{\genfrac{}{}{0pt}{}{D\left(\mathrm{𝟏}_n|\mathrm{𝟏}_m\right)=E|F}{DC}}{}}\mu (C,D)D`$
$`=`$ $`{\displaystyle \underset{E|FA|B}{}}\mu (A|B,E|F)E|F`$
$`+{\displaystyle \underset{E|FA|B}{}}{\displaystyle \underset{\genfrac{}{}{0pt}{}{D\left(\mathrm{𝟏}_n|\mathrm{𝟏}_m\right)=E|F}{DE|F}}{}}\left({\displaystyle \underset{\genfrac{}{}{0pt}{}{C\left(\mathrm{𝟏}_n|\mathrm{𝟏}_m\right)=A|B}{CD}}{}}\mu (C,D)\right)D.`$
The result follows then from the following equality
$$\underset{\genfrac{}{}{0pt}{}{C\left(\mathrm{𝟏}_n|\mathrm{𝟏}_m\right)=A|B}{CD}}{}\mu (C,D)=\underset{A|BCD}{}\mu (C,D)=0.$$
Induction and restriction define a graded product and coproduct on the space of $`G_0(𝚷,)=_{n0}G_0(\mathrm{𝕜}\mathrm{\Pi }_n,)`$. Define on the elements $`NG_0(\mathrm{𝕜}\mathrm{\Pi }_n,)`$ the operation
(3.4)
$$\mathrm{\Delta }(N)=\underset{k=0}{\overset{n}{}}\text{Res}_{k,nk}N$$
and for $`MG_0(\mathrm{𝕜}\mathrm{\Pi }_m,)`$,
(3.5)
$$NM=\text{Ind}_{n,m}NM.$$
$`G_0(𝚷,)`$ with the operation $`\mathrm{\Delta }`$ defines a coalgebra and $`G_0(𝚷,)`$ with the product of (3.5) defines an algebra structure. It is easily checked that the product and coproduct on $`G_0(𝚷,)`$ are not compatible as a bialgebra structure.
It is interesting to note that $`G_0(𝚷,)`$ endowed with the tensor product (3.2) and the coproduct $`\mathrm{\Delta }`$ from equation 3.4 does define a bialgebra. To highlight the relationship with $`\mathrm{𝖭𝖢𝖲𝗒𝗆}`$, we define an internal coproduct $`\mathrm{\Delta }^{}`$ on $`K_0(\mathrm{𝕜}\mathrm{\Pi }_n,)`$ which is the natural dual to equation (3.2). That is we define a map $`\mathrm{\Delta }^{}:K_0(\mathrm{𝕜}\mathrm{\Pi }_n,)K_0(\mathrm{𝕜}\mathrm{\Pi }_n,)K_0(\mathrm{𝕜}\mathrm{\Pi }_n,)`$ such that
(3.6)
$$\mathrm{\Delta }^{}(W_A)=\underset{BC=A}{}W_BW_C.$$
We can now show with the following theorem that $`K_0(𝚷,)`$ is a bialgebra and the simple modules in $`K_0(𝚷,)`$ correspond to the m-basis on $`\mathrm{𝖭𝖢𝖲𝗒𝗆}`$.
###### Theorem 3.5.
The ring $`K_0(𝚷,)`$ endowed with product $`MN:=\text{Ind}_{n,m}MN`$ and coproduct $`\mathrm{\Delta }^{}`$ of equation (3.6) defines a bialgebra. Moreover, the map $`F:K_0(𝚷,)\mathrm{𝖭𝖢𝖲𝗒𝗆}`$ given by $`F(W_A)=\text{m}_A`$ is an isomorphism of bialgebras.
###### Proof.
Recall that $`\mathrm{𝖭𝖢𝖲𝗒𝗆}`$ is a bialgebra linearly spanned by elements $`\text{m}_A`$ with the product defined by
$$\text{m}_A\text{m}_B=\underset{C(\mathrm{𝟏}_n|\mathrm{𝟏}_m)=A|B}{}\text{m}_C$$
and an inner coproduct defined by
$$\mathrm{\Delta }^{}(\text{m}_A)=\underset{BC=A}{}\text{m}_B\text{m}_C.$$
Equations (3.6) and (3.3) show that the map $`F(W_A)=\text{m}_A`$ is an isomorphism of bialgebras. ∎
This construction that we have presented here in the last two sections of defining an algebra from a lattice operation and looking at the modules is something that can be done in a more general setting and is a tool that can be used to analyze other Hopf algebras. This will be the subject of future work.
## 4. Existence of the $`𝐱_A`$ and Frobenius characteristic
We now prove our Theorem 1.3. It is useful at this point to introduce an intermediate basis of $`\mathrm{𝖭𝖢𝖲𝗒𝗆}`$. In , an analogue of the power sum basis is given by
(4.1)
$$𝐩_A=\underset{BA}{}\text{m}_B.$$
This basis has many nice properties.
###### Lemma 4.1.
The set $`\{𝐩_A:A\mathrm{\Pi }_n,n0\}`$ forms a basis of $`\mathrm{𝖭𝖢𝖲𝗒𝗆}`$ such that
$`\text{(i)}𝐩_A𝐩_B=𝐩_{A|B}`$
$`\text{(ii)}\mathrm{\Delta }^{}(𝐩_A)=𝐩_A𝐩_A.`$
###### Proof.
By triangularity, it is clear that the set forms a basis. Now, for $`A\mathrm{\Pi }_n`$ and $`B\mathrm{\Pi }_m`$ we have
$$𝐩_A𝐩_B=\underset{CA}{}\underset{DB}{}\text{m}_C\text{m}_D=\underset{CA}{}\underset{DB}{}\underset{E\mathbf{\hspace{0.17em}\hspace{0.17em}1}_n|\mathrm{𝟏}_m=C|D}{}\text{m}_E$$
Notice that we have that if $`E\mathrm{𝟏}_n|\mathrm{𝟏}_m=C|D`$, then $`EC|DA|B`$. Conversely, if $`EA|B`$, then we find unique $`C`$ and $`D`$ such that $`C|D=E\mathbf{\hspace{0.17em}1}_n|\mathrm{𝟏}_mA|B\mathbf{\hspace{0.17em}1}_n|\mathrm{𝟏}_m=(A\mathrm{𝟏}_n)|(B\mathrm{𝟏}_m)=A|B`$. This implies that the sum is equal to
$$𝐩_A𝐩_B=\underset{EA|B}{}\text{m}_E=𝐩_{A|B}.$$
For the second equality, we have
$`\mathrm{\Delta }^{}(𝐩_A)`$ $`={\displaystyle \underset{BA}{}}\mathrm{\Delta }^{}(\text{m}_B)={\displaystyle \underset{BA}{}}{\displaystyle \underset{CD=B}{}}\text{m}_C\text{m}_D`$
$`={\displaystyle \underset{CA}{}}{\displaystyle \underset{BA}{}}{\displaystyle \underset{D:CD=B}{}}\text{m}_C\text{m}_D={\displaystyle \underset{CA}{}}{\displaystyle \underset{DA}{}}\text{m}_C\text{m}_D`$
$`=𝐩_A𝐩_A`$
We finally define our basis. Let
(4.2)
$$𝐱_A=\underset{BA}{}\mu (B,A)𝐩_B.$$
By triangularity, the set $`\{𝐱_A:A\mathrm{\Pi }_n,n0\}`$ is an integral basis of $`\mathrm{𝖭𝖢𝖲𝗒𝗆}`$. We now see that this basis has the required properties.
###### Lemma 4.2.
$`\text{(i)}𝐱_A𝐱_B=𝐱_{A|B}`$
$`\text{(ii)}\mathrm{\Delta }^{}(𝐱_C)={\displaystyle \underset{AB=C}{}}𝐱_A𝐱_B.`$
###### Proof.
Using the same argument as in Lemma 2.5 we have
$`𝐱_A𝐱_B`$ $`={\displaystyle \underset{CA}{}}{\displaystyle \underset{DB}{}}\mu (C,A)\mu (D,B)𝐩_C𝐩_D`$
$`={\displaystyle \underset{CA}{}}{\displaystyle \underset{DB}{}}\mu (C,A)\mu (D,B)𝐩_{C|D}={\displaystyle \underset{EA|B}{}}\mu (E,A|B)𝐩_E=𝐱_{A|B}.`$
This shows the first identity. For the second, the left hand side of (ii) is
(4.3)
$$\mathrm{\Delta }^{}(𝐱_C)=\underset{EC}{}\mu (E,C)\mathrm{\Delta }^{}(𝐩_E)=\underset{EC}{}\mu (E,C)𝐩_E𝐩_E,$$
and the right hand side is
$$\underset{AB=C}{}𝐱_A𝐱_B=\underset{AB=C}{}\underset{\genfrac{}{}{0pt}{}{EA}{FB}}{}\mu (E,A)\mu (F,B)𝐩_E𝐩_F.$$
Let us isolate the coefficient of $`𝐩_E𝐩_F`$ in the sum above we get
(4.4) $`T_{E,F}^C`$ $`={\displaystyle \underset{\genfrac{}{}{0pt}{}{EAC}{FBC}}{}}{\displaystyle \underset{AB=C}{}}\mu (E,A)\mu (F,B)`$
$`={\displaystyle \underset{FBC}{}}\left({\displaystyle \underset{\genfrac{}{}{0pt}{}{EAC}{AB=C}}{}}\mu (E,A)\right)\mu (F,B).`$
By symmetry (interchanging the role of $`E`$ and $`F`$ if needed), we may assume that $`FE`$. In , Corollary 3.9.3 is dual to the following statement
$$\underset{\genfrac{}{}{0pt}{}{A\mathrm{𝟏}_n}{AB=\mathrm{𝟏}_n}}{}\mu (\mathrm{𝟎}_n,A)=\{\genfrac{}{}{0pt}{}{\mu (\mathrm{𝟎}_n,\mathrm{𝟏}_n)\text{if }B=\mathrm{𝟎}_n\text{,}}{0\text{otherwise.}}$$
where, as usual, $`\mathrm{𝟎}_n=\{1.2.\mathrm{}.n\}`$. This implies that the sum of in bracket in equation (4.4) is equal to
(4.5)
$$\underset{\genfrac{}{}{0pt}{}{EAC}{AB=C}}{}\mu (E,A)=\{\genfrac{}{}{0pt}{}{\mu (E,C)\text{if }B=E\text{,}}{0\text{otherwise.}}$$
This follows from the fact that $`\mu `$ is multiplicative and in general the interval $`[E,C]\mathrm{\Pi }_n`$ is isomorphic to a cartesian product of (smaller) partition lattices (see Example 3.9.4 in ). If we substitute this back in equation (4.4) we have two cases to consider. When $`FE`$, our assumption that $`FE`$ prohibits the possibility that $`FB=E`$. Thus we must always have $`BE`$ and in this case $`T_{E,F}^C=0`$. When $`F=E`$, the only value of $`B`$ where equation (4.5) does not vanish is when $`B=E=F`$ and we get $`T_{E,E}^C=\mu (E,C)\mu (E,E)=\mu (E,C)`$. If we compare this to equation (4.3) we conclude our proof of (ii). ∎
Notice that the character of the module (the trace of the matrix module) $`V_B`$ from formula (2.2) is given by the formula $`\chi ^{V_B}(A)=\delta _B(A)`$ when $`A\mathrm{𝕜}\mathrm{\Pi }_n`$ acts on $`V_B`$. We observe that equation (4.2) for $`𝐱_A`$ yields
$$𝐩_A=\underset{BA}{}𝐱_B=\underset{B}{}\chi ^{V_B}(A)𝐱_B.$$
This means that the characters for the simple modules for $`(𝚷,)`$ are encoded in the change of basis coefficients between the $`𝐩`$ and $`𝐱`$ basis.
Similarly, the character of the module $`W_B`$ when acted on by the element $`A\mathrm{𝕜}\mathrm{\Pi }_n`$ are given by the formula $`\chi ^{W_B}(A)=\delta _B(A)`$ from equation (3.1). Of course the defining relation of the $`𝐩`$ basis from equation (4.1) shows that
$$𝐩_A=\underset{B}{}\chi ^{W_B}(A)\text{m}_B.$$
We observe in this formula that the characters of the simple modules of $`(𝚷,)`$ are encoded in the change of basis coefficients between the $`𝐩`$ and m basis.
Both these formulas are in fairly close analogy with the formula for the expansion for the power basis in the Schur basis in the algebra of the symmetric functions. There the change of basis coefficients are the characters of the simple modules of the symmetric group. This shows that the $`𝐩`$-basis which was defined by Rosas and Sagan does represent the analogue of the power basis in the algebra of the symmetric functions and the $`𝐱`$ and the m bases encode in their coefficients the characters of the modules that they represent.
###### Remark 4.3.
One could also define a third algebra $`(\mathrm{𝕜}\mathrm{\Pi }_n,\mathrm{@})`$ where $`A\mathrm{@}B=\delta _{A=B}A`$ and construct the simple modules as we have done here for $`(\mathrm{𝕜}\mathrm{\Pi }_n,)`$ and $`(\mathrm{𝕜}\mathrm{\Pi }_n,)`$. This same construction shows that the simple modules of this algebra satisfy a tensor product, induction and restriction operations which make the Grothendieck ring (of the category of the finite dimensional projective modules) for this algebra isomorphic again to $`\mathrm{𝖭𝖢𝖲𝗒𝗆}`$ as a bialgebra where the simple modules behave as the elements $`𝐩_A\mathrm{𝖭𝖢𝖲𝗒𝗆}`$ and $`𝐩_A`$ is defined in (4.1).
###### Remark 4.4.
Summary of bases in $`\mathrm{𝖭𝖢𝖲𝗒𝗆}`$.
The $`𝐦`$ basis:
$`\text{m}_A\text{m}_B`$ $`=`$ $`{\displaystyle \underset{C(\mathrm{𝟏}_n|\mathrm{𝟏}_k)=A|B}{}}\text{m}_C`$
$`\mathrm{\Delta }(\text{m}_A)`$ $`=`$ $`{\displaystyle \underset{S[\mathrm{}(A)]}{}}\text{m}_{A_S}\text{m}_{A_{S^c}}`$
$`\mathrm{\Delta }^{}(\text{m}_A)`$ $`=`$ $`{\displaystyle \underset{BC=A}{}}\text{m}_B\text{m}_C`$
The $`𝐩`$ basis:
$`𝐩_A𝐩_B`$ $`=`$ $`𝐩_{A|B}`$
$`\mathrm{\Delta }(𝐩_A)`$ $`=`$ $`{\displaystyle \underset{S[\mathrm{}(A)]}{}}𝐩_{A_S}𝐩_{A_{S^c}}`$
$`\mathrm{\Delta }^{}(𝐩_A)`$ $`=`$ $`𝐩_A𝐩_A`$
The $`𝐱`$ basis:
$`𝐱_A𝐱_B`$ $`=`$ $`𝐱_{A|B}`$
$`\mathrm{\Delta }^{}(𝐱_A)`$ $`=`$ $`{\displaystyle \underset{BC=A}{}}𝐱_B𝐱_C`$
It would be interesting to find a formula for $`\mathrm{\Delta }(𝐱_A)`$. |
warning/0506/hep-th0506244.html | ar5iv | text | # Rotating Electric Classical Solutions of 2+1 D 𝑈(1) Einstein Maxwell Chern-Simons
## 1 Introduction
Several works have studied three dimensional classical gravitational configurations on topological and non topological field theories. The first works addressed Einstein theories, the well known AdS BTZ black hole , Einstein Maxwell Chern-Simons theory and rotating BTZ (see also ).
This work studies the classical solutions for a $`2+1`$D Einstein Maxwell Chern-Simons theory coupled to a gravitational massless scalar field (that is often interpreted as a dilaton field in string-frame). It therefore extends the work already done, both in Einstein Maxwell Chern-Simons theories , Einstein Maxwell theories , Einstein Maxwell theories with dilatonic potentials and the more recent Dilaton Einstein Maxwell theories .
Here we exclusively address pure electric solutions of $`3D`$ Einstein Maxwell Chern-Simons coupled to a massless scalar field. We try to present in a pedagogical way both general results and details of calculations. In particular our action resembles closely the one of the works together with a Chern-Simons term. However we start from a more generic action only particularizing the action due to the inexistence of other possible solutions.
The motivation to study our enlarged theory is two folded: the quantum consistence of the theory, and the embedding of a $`3D`$ system in a $`4D`$ world. First demanding quantum consistence of the theory we have to consider the Maxwell-Chern-Simons theory. Neither the pure Maxwell theory, neither the Chern-Simons theory are consistent at quantum level. If we start just with a Maxwell action, radiative (quantum) corrections will induce the Chern-Simons term and if we start with just a Chern-Simons action, quantum corrections will induce a Maxwell term, this correction is exact to all orders (see also for a review). Secondly our world is $`4D`$, therefore by counting degrees of freedom we need a gravitational scalar field in a $`3D`$ physical systems. Although several ways to embed $`2+1`$ dimensional systems in $`3+1`$ dimensions, the existence of a gravitational massless scalar field is rather well established. Considering a dimensional reduction scheme we obtain what is called Dilaton . Alternatively one can consider the gauging under some symmetry that effectively reduces the dimensionality of the problem, this is the example of the massless scalar field of the works on polarized cylindrical gravitational waves in $`3+1`$ gravity . It is not clear that this scalar field can always be interpreted as a dilaton field although for some particular actions it can be proved that it correspond to dilaton in string frame .
We also note that most of the literature in Abelian gauge Chern-Simons address (anti-)self-dual solutions. Here we address pure electric solutions.
The article is organized in the following way. In section 2 we present and discuss generic results of the Einstein Maxwell Chern-Simons theory coupled to a scalar field. First we introduce and justify the Action. From it we derive the equations of motion and choose a suitable metric parameterization. Also we derive the charge, angular momentum and the mass in the ADM formalism. In section 3 we solve the equations of motion in the Cartan-frame. In section 4 we compute the curvature, investigate the existence of singularities and horizons. Then in section 5 we compute the charge, angular momentum and mass for the configurations obtained. Finally in section 6 we summarize the solutions obtained and discuss them. In appendix A we introduce the Cartan Frame formalism (also known as non-coordinate frame) and derive the equations of motion and other useful formulae.
## 2 General Results
### 2.1 Action and EOM
We take a generic $`2+1`$D Einstein Gravity coupled to a massless scalar field with a Gauge Sector described by $`U(1)`$ Maxwell-Chern-Simons
$$\begin{array}{cc}\hfill S=& \frac{1}{2\pi }_Md^3x\{\sqrt{g}[e^{a\varphi }(R+2\lambda (\varphi )^2)e^{b\varphi }\mathrm{\Lambda }\hfill \\ & +\widehat{ϵ}\frac{e^{c\varphi }}{2}(F_{\mu \nu }F^{\mu \nu }+J^\mu A_\mu )]\widehat{ϵ}\frac{m}{2}ϵ^{\mu \nu \lambda }A_\mu F_{\nu \lambda }\}\hfill \end{array}$$
(2.1)
where $`a`$, $`b`$, $`c`$, $`\lambda `$ and the cosmological constant $`\mathrm{\Lambda }`$ are numerical parameters of the theory. $`\widehat{ϵ}=\pm 1`$ simply sets the relative sign between the gauge sector and the gravitational sector.
Varying this action in relation to the fields $`A_\mu `$, $`g^{\mu \nu }`$ and $`\varphi `$ we obtain the equations of motion, i.e. the Maxwell, Einstein and scalar field equations
$$\begin{array}{ccc}\hfill _\alpha (\sqrt{g}e^{c\varphi }F^{\alpha \mu })+\frac{m}{2}ϵ^{\mu \alpha \beta }F_{\alpha \beta }& =& \sqrt{g}e^{c\varphi }J^\mu \hfill \\ \hfill G_{\mu \nu }a_\mu _\nu \varphi +ag_{\mu \nu }^2\varphi +(\lambda a^2)_\mu \varphi _\nu \varphi & & \\ \hfill \left(\frac{\lambda }{2}a^2\right)g_{\mu \nu }(\varphi )^2+\frac{1}{2}e^{(ba)\varphi }g_{\mu \nu }\mathrm{\Lambda }& =& 2e^{(ca)\varphi }T_{\mu \nu }\hfill \\ \hfill e^{a\varphi }\left[2(2a^2\lambda )^2\varphi +2a(2a^2\lambda )(\varphi )^2\right]+(3ab)e^{b\varphi }\mathrm{\Lambda }& =& \widehat{ϵ}(a+c)e^{c\varphi }F^2\hfill \end{array}$$
(2.2)
where the Einstein and Stress-Energy tensors are defined as
$$\begin{array}{ccc}\hfill G_{\mu \nu }& =& R_{\mu \nu }\frac{1}{2}g_{\mu \nu }R\hfill \\ \hfill T_{\mu \nu }& =& \widehat{ϵ}\left(F_{\mu \alpha }F_\nu ^\alpha \frac{1}{4}g_{\mu \nu }F^2\right)\hfill \end{array}$$
(2.3)
and the covariant derivative and Laplacian are as usual
$$\begin{array}{ccc}\hfill _\mu _\nu \varphi & =& _\mu _\nu \varphi \mathrm{\Gamma }_{\mu \nu }^\alpha _\alpha \varphi \hfill \\ \hfill ^2\varphi & =& _\alpha ^\alpha \varphi +\mathrm{\Gamma }_{\alpha \beta }^\alpha ^\beta \varphi \hfill \end{array}$$
(2.4)
Note that the scalar field equations presented are obtained from the usual variation of the action with respect to $`\varphi `$
$$e^{a\varphi }\left[aR2\lambda ^2\varphi a\lambda (\varphi )^2\right]be^{b\varphi }\mathrm{\Lambda }=\widehat{ϵ}ce^{c\varphi }F^2$$
(2.5)
summed with the contraction of the 3 Einstein equations with the metric times $`2a`$. In this way the gravitational curvature is not present in equation (2.2).
Our convention for the Ricci tensor is
$$R_{\mu \nu }=\mathrm{\Gamma }_{\mu \alpha ,\nu }^\alpha +\mathrm{\Gamma }_{\mu \nu ,\alpha }^\alpha \mathrm{\Gamma }_{\mu \beta }^\alpha \mathrm{\Gamma }_{\nu \alpha }^\beta +\mathrm{\Gamma }_{\beta \alpha }^\alpha \mathrm{\Gamma }_{\mu \nu }^\beta $$
(2.6)
we note that when considering a cosmological constant $`\mathrm{\Lambda }`$ the symmetric definition of the Ricci tensor, maintaining the same metric signature, is not equivalent and will account for the opposite sign of $`\mathrm{\Lambda }`$. In order to justify this choice, in the next subsection, we give the example of $`3`$-dimensional deSitter space, a known and well studied example with $`\mathrm{\Lambda }>0`$.
### 2.2 Metric, Ricci Tensor and Maxwell Tensor
We take several parameterizations of a radial symmetric metric, in polar coordinates $`x^0=t`$, $`x^1=r`$ and $`x^2=\phi `$ of the form
$$ds^2=g_{tt}dt^2+dr^2+g_{\phi \phi }d\phi ^2+2g_{t\phi }dtd\phi $$
(2.7)
with signature $`(,+,+)`$.
The Antisymmetric tensor has only the non vanishing components
$$F_{tr}=E_{}F_{r\phi }=B_{}$$
(2.8)
All the functions $`g_{tt}`$, $`g_{\phi \phi }`$, $`g_{t\phi }`$, $`E_{}`$, $`B_{}`$ and $`\varphi `$ are radial symmetric, i.e. are $`r`$ dependent only.
There is a couple of important well establish points to stress to fully justify this ansatz.
The motivation of introducing the $`g_{t\phi }`$ component of the metric is due to the Maxwell equations, in the presence of the Chern-Simons term (without external currents), not allowing for solutions $`B_{}=0`$ or $`E_{}=0`$ when $`g_{t\phi }=0`$ (both must be null or both must be present). So when there is a Chern-Simons term in the action and we are considering only Electric or only Magnetic fields, we must have $`g_{t\phi }0`$, otherwise both fields are null. In physical terms means that the space-time is rotating, although it can still be stationary as long as $`g_{t\phi }`$ does not depend on the time coordinate.
Also one may consider a non null $`F_{t\phi }`$ but for the metric parameterizations considered here the Maxwell Equation in (2.2) for $`\mu =1`$ imposes it to be null.
Finally it is important to stress that one can add a generic parameterization for $`g_{rr}=1/L^2`$ by introducing a new radial coordinate $`\rho `$ such that $`d\rho /dr=L`$. This accounts for a choice of coordinates and therefore does not change the physical results presented here.
Although in $`4D`$ space-time the choice of metric (most positive or most negative diagonal) is not relevant, in $`3D`$ space-time one needs extra care in the relative definitions between the metric and remaining tensor fields. The reader may also note that depending on the choice of $`3D`$ Minkowski metric the determinant is positive (for most negative diagonal) or negative (for most positive diagonal). In (2.7) we choose the last case to maintain the determinant of the metric negative. To justify the choice of the Ricci tensor (2.6) and clear any confusions concerning its definition we present a simple pedagogical example of the well known dS geometry which has positive cosmological constant. We consider a cosmological Einstein action
$$S_E=d^3x\sqrt{g}\left(R2\mathrm{\Lambda }\right)$$
(2.9)
and a dS metric for an observer at r=0 corresponding to a cosmological constant $`\mathrm{\Lambda }=+1`$, of the form
$$ds^2=(1r^2)dt^2+\frac{1}{(1r^2)}dr^2+r^2d\phi ^2$$
(2.10)
with signature $`(,+,+)`$ near the origin (where the observer is) and determinant $`|g|=r^2`$. Varying the action with respect to $`g^{\mu \nu }`$ we obtain the well know equations of motion
$$G_{\mu \nu }+\mathrm{\Lambda }g_{\mu \nu }=0$$
(2.11)
where $`G_{\mu \nu }=R_{\mu \nu }g_{\mu \nu }R/2`$ is the usual Einstein tensor. For the given metric, computing explicitly the einstein tensor, we obtain $`G_{00}=1r^2`$, $`G_{11}=1/(1r^2)`$ and $`G_{22}=r^2`$. This reads
$$G_{\mu \nu }=g_{\mu \nu }$$
(2.12)
Therefore the cosmological constant is uniquely define trough the equations of motion as $`\mathrm{\Lambda }=+1`$. Maintaining the metric signature and the action and considering the symmetric definition of the Ricci tensor $`\stackrel{~}{R}_{\mu \nu }=R_{\mu \nu }`$ we would obtain $`\stackrel{~}{G}_{\mu \nu }=G_{\mu \nu }`$ and hence $`\mathrm{\Lambda }=1`$. Together with the definition $`\stackrel{~}{R}`$, if we swap the signature of the metric to $`(+,,)`$ maintaining the action or if we maintain the signature of the metric and change the action to $`\stackrel{~}{S}_E=(R+2\mathrm{\Lambda })`$ we would obtain $`\mathrm{\Lambda }=+1`$. Also using the definition $`R`$, swapping the signature of the metric to $`(+,,)`$ and considering the action $`\stackrel{~}{S}_E`$ we would obtain $`\mathrm{\Lambda }=+1`$.
So we conclude that the choices of the definition of the Ricci tensor, the metric signature and the relative sign of the cosmological constant and the gravitational curvature in the action are not all equivalent. Resuming, we choose the definition of the Ricci tensor given by (2.6), the metric signature $`(,+,+)`$ and an action of the form (2.1).
Finally we briefly discuss the relative sign between the several terms in the action. First we note that we consider opposite signs between the Chern-Simons term the Maxwell term. This is to ensure that the photon mass is real $`(^2m^2)F^{}=0`$ , if they have the same sign we would obtain imaginary (tachyonic) masses. In particular this choice sets the sign of the angular momentum $`J`$, as we will se our solutions have $`Jm`$ (or $`\mathrm{sign}(m)`$). This is an effect of parity violation and is expected because the Chern-Simons term violates parity in the gauge sector. If we change the relative sign between the Maxwell term and the Chern-Simons term the only effect on the solutions is to change the sign of the angular momentum. However as we explained this accounts for the photon to become a tachyon, for this reason we fixed this choice.
$`\widehat{ϵ}=\pm 1`$ sets the relative sign between the gauge sector (Maxwell term $`F^2`$) and the Einstein term ($`R`$). Choosing $`\widehat{ϵ}=+1`$ or $`\widehat{ϵ}=1`$ does not change the expressions for the solutions, nevertheless the validity range for the parameters will change significantly, therefore the physical interpretation of the results as well. Also it is interesting to note that upon quantization the sign of the Maxwell is relevant. If we have $`\widehat{ϵ}=1`$ we obtain the standard Hamiltonian and excited states of the gauge fields will have positive energy for Bose-Einstein spin-statistics, while for $`\widehat{ϵ}=+1`$ the excited states for the gauge fields will only hold positive energy for Fermi-Dirac spin-statistics. In this case the gauge fields have the wrong spin-statistics and for that reason are commonly called ghosts. It is quite interesting that different choice of signs will also at classical level hold significant differences as we will see in detail.
### 2.3 Mass, Charge and Angular Momentum
We are going to use the ADM formalism (see ), so we rewrite the line element using a generic parameterization<sup>1</sup><sup>1</sup>1This metric parameterization is not unique but it accounts for the most generic parameterization for a stationary radial symmetric $`2+1D`$ metric
$$ds^2=f^2dt^2+dr^2+h^2(d\phi +Adt)^2$$
(2.13)
and considering the Hamiltonian form of the action
$$S=2\pi \mathrm{\Delta }t𝑑r\left[f+A^\phi +A_0𝒢\right]+S_{}$$
(2.14)
where $`S_{}`$ stands for boundary terms due to the integration by parts of the terms containing $`f^{}`$, $`f^{\prime \prime }`$, $`A^{}`$ and $`A_0^{}`$
$$S_{}=\frac{1}{2\pi }_Md^2x\left[f\mathrm{\hspace{0.17em}2}e^{a\varphi }(2ah\varphi ^{}+h^{})+Ae^{a\varphi }\mathrm{\Pi }_G+\widehat{ϵ}A_0\left(\mathrm{\Pi }_{EM}\frac{m}{2}A_\phi \right)\right]$$
(2.15)
and the Hamiltonian, Momentum and Gauss constraints are respectively
$$\begin{array}{ccc}\hfill & =& \frac{2\mathrm{\Pi }_G^2}{h^3}e^{a\varphi }2a\left(h\varphi ^{}e^{a\varphi }\right)^{}2h^{\prime \prime }e^{a\varphi }+2\lambda h(\varphi ^{})^2e^{a\varphi }+\mathrm{\Lambda }he^{b\varphi }\hfill \\ & +& \widehat{ϵ}\left(\frac{e^{c\varphi }}{h}\left(\mathrm{\Pi }_{EM}+\frac{m}{2}A_\phi \right)^2+he^{c\varphi }(A_\phi ^{})^2\right)\hfill \\ \hfill ^\phi & =& (\mathrm{\Pi }_Ge^{a\varphi })^{}\widehat{ϵ}\left(\mathrm{\Pi }_{EM}+\frac{m}{2}A_\phi \right)A_\phi ^{}\hfill \\ \hfill 𝒢& =& \widehat{ϵ}\left(\mathrm{\Pi }_{EM}\frac{m}{2}A_\phi \right)^{}.\hfill \end{array}$$
(2.16)
$`\sqrt{g}=hf`$ and the induced $`2D`$ metric is simply $`h_{ij}=\mathrm{diag}(1,h^2)`$. The prime () means the usual derivation ($`_r`$) with respect to $`r`$. We note that $`\widehat{ϵ}`$ in the gauss constraint is optional once it is a constraint of the gauge sectors only.
For the rotating radially symmetric configurations considered in this work (see subsections 2.1 and 2.2) the only non vanishing gravitational canonical momenta conjugate to $`h_{ij}`$ is $`\pi _G^{r\phi }`$ (conjugate to $`h_{r\phi }`$) such that
$$\mathrm{\Pi }_G=\mathrm{Tr}(\pi _G)=(\pi _G)_\phi ^r=\frac{h^3A^{}}{f}$$
(2.17)
and the only non vanishing gauge canonical momenta conjugate to $`A_i`$ is $`\pi _{EM}^r=\delta S/\delta (_0A_r)`$ (conjugate to $`A_r`$) such that
$$\mathrm{\Pi }_{EM}=\widehat{ϵ}\left(\frac{m}{2}A_\phi \right)$$
(2.18)
The contravariant Electric and Magnetic densities are defined as
$$\begin{array}{ccc}\hfill & =& hfe^{c\varphi }F^{0r}=\widehat{ϵ}\left(\mathrm{\Pi }_{EM}+\frac{m}{2}A_\phi \right)\hfill \\ \hfill & =& hfe^{c\varphi }ϵ^{r\phi }F_{r\phi }=hfe^{c\varphi }A_\phi ^{}\hfill \end{array}$$
(2.19)
For completeness we also note that the contravariant current densities are defined as
$$𝒥^\mu =hfe^{c\varphi }J^\mu $$
(2.20)
There is a couple of important points that should be stressed. Generally, due to the rotation, magnetic configurations generate a magnetic field and magnetic configurations generate an electric field. However we will solve our equations in the Cartan frame such that for given fields $`E`$ and $`B`$ in the Cartan frame we obtain $`=he^{c\varphi }E`$ and $`=h^2fe^{c\varphi }B`$. Therefore we don’t actually have mixing between electric and magnetic fields (see appendix A). Ww also note that in $`3D`$ the magnetic field is a scalar that corresponds in $`4D`$ to the $`z`$-component of the magnetic field, this means the magnetic field perpendicular to the $`2D`$ spatial coordinates. In our configurations it is null.
The generic gauge canonical momenta are $`\pi _{EM}^i=\widehat{ϵ}(hfe^{c\varphi }F^{0i}mϵ^{ij}A_j/2)`$ and therefore $`\pi _{EM}^\phi `$ is not generally null. However we are only studying configurations in which $`F_{t\phi }=_tA_\phi _\phi A_t=0`$ (see discussion on subsection 2.2) such that $`\pi _{EM}^\phi =m\widehat{ϵ}A_r`$. Since we are considering only rotating radial symmetric configurations we consider that all the gauge fields are radial functions, furthermore we still have a radial gauge freedom, this means that a gauge transformation $`\mathrm{\Lambda }(r)`$ depending on the radius only has the effect $`A_rA_r+\mathrm{\Lambda }^{}(r)`$ and does not change any of the physical quantities. Therefore we can without lost of generality gauge fix $`\pi _{EM}^\phi =A_r=0`$.
As a final remark note that in the pure Maxwell theory ($`m=0`$) the canonical momentum is proportional to the Electric density $`\mathrm{\Pi }_{\mathrm{Maxwell}}=\widehat{ϵ}`$ such that this density is itself a canonical variable, with the Chern-Simons term this is no longer true.
Varying both the action $`S`$ and the boundary action $`A_{}`$ with respect to the canonical dynamical variables ($`h,\mathrm{\Pi }_G,\varphi ,\mathrm{\Pi }_{EM},A_\phi `$) one obtains a boundary variation
$$\delta S_{}=2\pi \mathrm{\Delta }t(f\delta M+A_0\delta Q+A\delta J)$$
(2.21)
where $`M`$, $`Q`$ and $`J`$ are the Mass, Charge and Angular Momentum of the configuration and $``$ stands for the one-dimensional spatial boundary of the spatial manifold. Their variation is
$$\begin{array}{ccc}\hfill \delta M& =& 2\delta \left(he^{a\varphi }\right)^{}+4\lambda h\varphi ^{}e^{a\varphi }\delta \varphi +2\widehat{ϵ}he^{c\varphi }A_\phi ^{}\delta A_\phi |_{}\hfill \\ \hfill \delta Q& =& 2\widehat{ϵ}\delta \left(\mathrm{\Pi }_{EM}\frac{m}{2}A_\phi \right)|_{}\hfill \\ \hfill \delta J& =& 2\delta \left(\mathrm{\Pi }_Ge^{a\varphi }\right)2\widehat{ϵ}\left(\mathrm{\Pi }_{EM}+\frac{m}{2}A_\phi \right)\delta A_\phi |_{}\hfill \end{array}$$
(2.22)
In order to exist well defined classical minimum it is necessary that these variations vanish. We need either to add a boundary action that cancels these variations or to demand them (the variations) to vanish at the boundary. The later is usually a very strong condition and accounts for having expressions for M, Q and J to be constants (meaning $`r`$ independent). In the absence of external currents the charge $`Q`$ is necessarily a constant since the Gauss’ law is expressed as a total derivative. Accounting with the charge expression, the angular momentum $`J`$ is also expressed as a total derivative and is therefore a constant as well. For the case of the mass $`M`$ this is no longer true and we need to add a suitable boundary action. In the presence of external currents neither $`Q`$ nor $`J`$ are generally constants since the Gauss’ law includes the external charge and is no longer a total derivative, here we are not addressing this case.
Considering the above procedure we obtain
$$\begin{array}{ccc}\hfill M& =& 2(he^{a\varphi })^{}+4\lambda h\varphi \varphi ^{}e^{a\varphi }+2\widehat{ϵ}he^{c\varphi }A_\phi A_\phi ^{}|_{r0}^r\mathrm{}\hfill \\ \hfill Q& =& 2\widehat{ϵ}\left(\mathrm{\Pi }_{EM}\frac{m}{2}A_\phi \right)|_{r0}^r\mathrm{}\hfill \\ \hfill J& =& 2\mathrm{\Pi }_Ge^{a\varphi }2\widehat{ϵ}\left(Q+\frac{m}{2}A_\phi \right)A_\phi |_{r0}^r\mathrm{}\hfill \end{array}$$
(2.23)
where we used the fact that once the charge constraint in equation (2.22) is taken care, the charge variation vanishes $`\delta Q=0`$, and used the expression for the charge to replace $`\mathrm{\Pi }_{EM}=Q+mA_\phi /2`$ in the second term of the equation for the angular momentum variation in order to get a variation of $`A_\phi `$ only. We are considering two disconnected boundaries, the spatial infinite $`r=\mathrm{}`$ and the singularity at the origin $`r=0`$. We note that these two boundaries have opposite orientations, such that their contributions add up.
As for the mass expression we have to be careful with what fields are fixed and what fields vary upon a functional variation. The correct expression should be
$$M=2(he^{a\varphi })^{}+4\lambda h\varphi \widehat{\varphi }^{}e^{a\widehat{\varphi }}+2\widehat{ϵ}he^{c\widehat{\varphi }}A_\phi \widehat{A}_\phi ^{}|_{r0}^r\mathrm{}$$
(2.24)
where the hatted fields are fixed at the two boundaries ($`r0`$ and $`r\mathrm{}`$), i.e. upon a functional variation of the mass we obtain the correct expression (2.22).
### 2.4 Geodesics and Horizons
To compute the geodesics we use the variational principle presented in , so we consider the constant functional
$$K=g_{\mu \nu }\frac{x^\mu }{ds}\frac{x^\nu }{ds}=\kappa =\{\begin{array}{cc}\hfill 0& \mathrm{for}\mathrm{lightlike}(\mathrm{null})\mathrm{geodesics}\hfill \\ \hfill 1& \mathrm{for}\mathrm{timelike}\mathrm{geodesics}\hfill \\ \hfill +1& \mathrm{for}\mathrm{spacelike}\mathrm{geodesics}\hfill \end{array}$$
(2.25)
where the derivatives are with respect to a affine parameter $`s`$. We minimize $`K`$ solving the Euler-Lagrange equations $`\frac{\delta K}{\delta x^\mu }\frac{d}{ds}\left(\frac{\delta K}{\delta \dot{x}^\mu }\right)=0`$. Since our solutions are both cylindrically symmetric and stationary (only depend on $`r`$, the radial coordinate) we have that the equations for $`\mu =t,\phi `$ lead respectively to the first integrals of motion
$$\{\begin{array}{ccc}\hfill g_{00}\frac{dt}{ds}+g_{02}\frac{d\phi }{ds}& =& E\hfill \\ \hfill g_{22}\frac{d\phi }{ds}+g_{02}\frac{dt}{ds}& =& L\hfill \end{array}\{\begin{array}{ccc}\hfill \frac{dt}{ds}& =& \frac{Eg_{22}Lg_{02}}{g}\hfill \\ \hfill \frac{d\phi }{ds}& =& \frac{Lg_{00}Eg_{02}}{g}\hfill \end{array}$$
(2.26)
with $`2E=p_t=\frac{\delta K}{\delta \dot{t}}`$ and $`2L=p_\phi =\frac{\delta K}{\delta \dot{\phi }}`$ being constants of motion, the energy and angular momentum respectively (here we rescaled them by a factor of $`2`$ in order to simplify the expressions). Using the two equations (2.26) in (2.25) we obtain an expression for $`dr/ds`$
$$\left(\frac{dr}{ds}\right)^2=k\frac{L^2g_{00}2ELg_{02}+E^2g_{22}}{g}$$
(2.27)
being $`g`$ the determinant of the metric $`g=g_{00}g_{22}g_{02}^2`$.
Since we are looking for stationary polar symmetric solutions $`\dot{t}`$ and $`\dot{\phi }`$ can be expressed in terms of the radial variable $`r`$ only, $`(d/ds)/(dr/ds)=d/dr`$. From the equations for $`t`$ and $`\phi `$ (2.26) we obtain the differential equations
$$\begin{array}{ccc}\hfill t^{}(r)& =& \pm \frac{Eg_{22}Lg_{02}}{\sqrt{g\left(g\kappa L^2g_{00}+2ELg_{02}E^2g_{22}\right)}}\hfill \\ \hfill \phi ^{}(r)& =& \pm \frac{Lg_{00}Eg_{02}}{\sqrt{g\left(g\kappa L^2g_{00}+2ELg_{02}E^2g_{22}\right)}}\hfill \end{array}$$
(2.28)
Solving these equation one obtains the $`t`$ and $`\phi `$ dependence on $`r`$. We can also compute the radial velocity $`\dot{r}=(dr/ds)/(dt/ds)`$ and angular velocity $`\dot{\phi }=(d\phi /ds)/(dt/ds)`$
$$\begin{array}{ccc}\hfill \dot{r}(r)& =& \pm \frac{\sqrt{g\left(g\kappa L^2g_{00}+2ELg_{02}E^2g_{22}\right)}}{Eg_{22}Lg_{02}}\hfill \\ \hfill \dot{\phi }(r)& =& \frac{Lg_{00}Eg_{02}}{Eg_{22}Lg_{02}}\hfill \end{array}$$
(2.29)
We note that these solutions are for an external observer (at rest far away from the singularity). Then the first equation is particular useful, when $`\dot{r}=0`$ we are in the presence either of a turning point on the trajectory, or of a horizon (in which case the geodesics at the rest frame of the travelling observer hits the singularity). We also note that at the singularity, if $`\dot{r}`$ is null the singularity is not naked, meaning that an external observer sees the particle stopping when arriving to the singularity. While if $`\dot{r}`$ has some positive value at the singularity we have a naked singularity since an external observer can actually see it without reaching it. We are using this results to inquire if we have an horizon or not.
## 3 Electric Solutions
Here we will look for pure Electric solution without external currents, hence we set $`B=B_{}`$, being $`B`$ the magnetic field in the Cartan frame and $`B_{}`$ the magnetic field in the original frame. We will be working in the Cartan frame and at the end of each subsection we will summarize our results in the original frame. The equations of motion in the Cartan frame are computed in appendix A and are equivalent to the equations of motion as presented in subsection 2.1.
From the first Maxwell Equation (A.24) we obtain that
$$\gamma =me^{c\varphi }$$
(3.1)
Using (3.1) in (A.26) one gets that $`\beta =c\varphi ^{}/2`$ and from the definition of $`\beta `$ (see (A.21) in appendix) we get the solution for $`h`$
$$h=c_he^{\frac{c}{2}\varphi }$$
(3.2)
where $`c_h`$ is a free integration constant.
Now we get from the second Maxwell Equation (A.25) that
$$E=\chi e^{\frac{3}{2}c\varphi }$$
(3.3)
where $`\chi `$ is an integration constant. Note that without loss of generality we included $`c_h`$ in the definition of this constant. There is a very important conclusion to take from this last equation, trivial solutions for the scalar field ($`\varphi =\mathrm{constant}`$) holds in the Cartan frame a uniform (constant) electric field $`E`$ in all space, this conclusion was firstly obtained in . Although for completeness we address trivial solutions we will first address non-trivial solutions for the scalar field which is the main objective of this work.
### 3.1 Non-Trivial Solutions for the Scalar Field
We will now address the full equations considering the generic equations. The three Einstein (A.27-A.29) and scalar field equations (A.30) read now
$`(a+{\displaystyle \frac{c}{2}})\varphi ^{\prime \prime }+(a^2{\displaystyle \frac{\lambda }{2}}+{\displaystyle \frac{c^2}{4}})(\varphi ^{})^2+{\displaystyle \frac{m^2}{4}}e^{2c\varphi }+{\displaystyle \frac{\mathrm{\Lambda }}{2}}e^{(ba)\varphi }`$ $`=`$ $`\widehat{ϵ}\chi ^2e^{(a2c)\varphi }`$ (3.4)
$`a\varphi ^{\prime \prime }+(a^2{\displaystyle \frac{\lambda }{2}})(\varphi ^{})^2+\alpha ^2+\alpha ^{}{\displaystyle \frac{3m^2}{4}}e^{2c\varphi }+{\displaystyle \frac{\mathrm{\Lambda }}{2}}e^{(ba)\varphi }`$ $`=`$ $`\widehat{ϵ}\chi ^2e^{(a2c)\varphi }`$ (3.5)
$`{\displaystyle \frac{\lambda }{2}}(\varphi ^{})^2+{\displaystyle \frac{c}{2}}\alpha \varphi ^{}+{\displaystyle \frac{m^2}{4}}e^{2c\varphi }+{\displaystyle \frac{\mathrm{\Lambda }}{2}}e^{(ba)\varphi }`$ $`=`$ $`\widehat{ϵ}\chi ^2e^{(a2c)\varphi }`$ (3.6)
$`(4a^2\lambda )\varphi ^{\prime \prime }+a(4a^22\lambda )(\varphi ^{})^2+(3ab)\mathrm{\Lambda }e^{(ba)\varphi }`$ $`=`$ $`\widehat{ϵ}(a+c)\chi ^2e^{(a2c)\varphi }`$ (3.7)
The main problem to solve these equations is to make them compatible with each other in order to give a non trivial solution. For $`a=b=c`$, for $`a=0`$ (any $`b`$ and $`c`$), for $`b=0`$ (any $`a`$ and $`c`$) and $`c=0`$ (any $`a`$ and $`b`$) these equations hold that the scalar field has only trivial solutions, i.e. it must be a constant. Trivial solutions will be addressed in the next subsections. For the particular cases $`c=0`$ with $`a=b`$ and $`a=b=2c`$ solutions do exist but hold that the scalar field is purely imaginary.
The better way to properly understood the structure of the equations is the following. The third equation (3.6) can be algebraically solved in $`\alpha `$ which solution is then plugged into the second equation (3.5). Then to obtain a solution for the $`\varphi `$ we can make a linear combination of the remaining three equations obtaining a simpler equation. The main problem then is to ensure that the solution is compatible with the original equations (or equivalently with different linear combinations of the original equations). This procedure gives very few choices for non-trivial solutions.
We only found non-trivial solutions for the case
$$\begin{array}{ccc}\hfill a& =& 0\hfill \\ \hfill c& =& \frac{b}{2}\hfill \\ \hfill \lambda & & \frac{b^2}{8}\hfill \end{array}$$
(3.8)
For $`b^2=8\lambda `$ does not exist a non-trivial solution either. We note that for the choice of equation (3.8) we are not working with dilaton Einstein theory. Our action is more similar to what is commonly know as a gravitational scalar field and the cosmological constant term resembles a Dilaton potential <sup>2</sup><sup>2</sup>2Thanks to Dmitri Gal’tsov for this remark.
Given this ansatz we combine (3.4) with (3.7) obtaining
$$\varphi ^{}=\pm \sqrt{c_1}e^{b\varphi }$$
(3.9)
such that
$$\varphi =\frac{2}{b}\mathrm{ln}(c_\varphi (rr_0))$$
(3.10)
Here
$$c_\varphi =\frac{|b|}{2}\sqrt{c_1}c_1=2\frac{b^2(\widehat{ϵ}\chi ^2+2\mathrm{\Lambda })+2\lambda (4\widehat{ϵ}\chi ^2+2\mathrm{\Lambda }+m^2)}{\lambda (b^28\lambda )}$$
(3.11)
and without loss of generality we set the integration constant $`r_0=0`$ since it represents only a shift in the radial coordinate and all the solutions depend on the $`\varphi `$ exponentials. Note that the choice of sign in (3.9) depends on the sign of $`b`$ such that in (3.10) the argument of the logarithm is positive. Also we have to ensure that $`c_1`$ is positive defined. Before doing so we use the $`\varphi `$ solution (3.10) in (3.4). In order the equation to be solved we have to impose
$$\chi ^2=\widehat{ϵ}\frac{2\mathrm{\Lambda }(b^2+12\lambda )+4\lambda m^2}{b^2+24\lambda }$$
(3.12)
Now $`c_1`$ becomes
$$c_1=4\frac{m^26\mathrm{\Lambda }}{b^2+24\lambda }$$
(3.13)
From (3.6) and the definition $`\gamma =A^{}h/f`$ (A.21) we get that
$$\alpha =\left(16\frac{\lambda }{b^2}+1\right)\frac{1}{2r}$$
(3.14)
Therefore from the definition of $`\alpha =f^{}/f`$ (see (A.21) in the appendix) we obtain the solution for $`f`$
$$f=c_fr^{\frac{8\lambda }{b^2}\frac{1}{2}}$$
(3.15)
from (3.2) we obtain the solution for $`h`$
$$h=c_h\sqrt{r}$$
(3.16)
and from (3.1) we get the solution for $`A`$
$$A=c_Ar^{\frac{8\lambda }{b^2}1}+c_A_{\mathrm{}}$$
(3.17)
where
$$c_A=\frac{mc_f}{c_h\left(\frac{8\lambda }{b^2}1\right)}\sqrt{\frac{1+\frac{24\lambda }{b^2}}{m^26\mathrm{\Lambda }}}$$
(3.18)
$`c_f`$, $`c_h`$ and $`c_A_{\mathrm{}}`$ are free constants.
Replacing these solutions in the remaining equation (3.5) and demanding it to be obeyed we get that
$$\lambda _\pm =\frac{b^2}{8}\frac{3\mathrm{\Lambda }\sqrt{\mathrm{\Lambda }(2m^23\mathrm{\Lambda })}}{m^26\mathrm{\Lambda }}$$
(3.19)
We have to ensure that all these relations are possible and that do not correspond to trivial solutions, in particular that $`\chi ^2>0`$ and $`C_1>0`$. Therefore for each $`\widehat{ϵ}=\pm 1`$ we have to choose the solution $`\lambda _{\widehat{ϵ}}`$ getting
$$\begin{array}{ccc}\hfill \chi ^2& =& \frac{1}{2}\left[\widehat{ϵ}\mathrm{\Lambda }+\sqrt{\mathrm{\Lambda }(2m^23\mathrm{\Lambda })}\right]\hfill \\ \hfill C_1& =& \frac{4}{b^2}\left[3\mathrm{\Lambda }+m^2+3\widehat{ϵ}\sqrt{\mathrm{\Lambda }(2m^23\mathrm{\Lambda })}\right]\hfill \end{array}$$
(3.20)
Demanding positiveness of these expressions hold, independently of $`\widehat{ϵ}`$ the same constraint on the cosmological constant $`\mathrm{\Lambda }`$ and topological mass $`m`$
$$0<\mathrm{\Lambda }<\frac{m^2}{2}$$
(3.21)
For the particular value of $`\mathrm{\Lambda }=m^2/6`$ some of the expressions previously computed are not well defined. It is necessary to rederive the solution using the same method. For $`\widehat{ϵ}=+1`$ we obtain
$$\mathrm{\Lambda }=m^2/6C_1=\frac{12m^2}{b^2}\chi ^2=\frac{m^2}{6}\lambda =\frac{b^2}{24}C_A=\frac{\sqrt{3}c_f}{c_h}.$$
(3.22)
All the other solutions remain the same up to replacement of the above constants. For $`\widehat{ϵ}=1`$ there are no allowed solutions at $`\mathrm{\Lambda }=m^2/6`$.
For convenience we define the parameter $`p`$ which depends only on the ratio $`\mathrm{\Lambda }/m^2`$
$$p=8\frac{\lambda }{b^2}=\frac{3x+\widehat{ϵ}\sqrt{x(23x)}}{16x}x=\frac{\mathrm{\Lambda }}{m^2}$$
(3.23)
For clarity we summarize and rewrite the solutions computed above in the original frame,
$$\begin{array}{ccc}\hfill \varphi & =& \frac{2}{b}\mathrm{ln}(C_\varphi r)\hfill \\ \hfill h& =& C_h\sqrt{r}\hfill \\ \hfill f& =& C_fr^{p\frac{1}{2}}\hfill \\ \hfill A& =& C_Ar^{p1}+\theta \hfill \\ \hfill E_{}& =& C_Er^{p2}\hfill \\ \hfill A_0& =& \frac{C_E}{p1}r^{p1}\hfill \end{array}$$
(3.24)
where for convenience we rename the variables and integration constants. $`C_h`$, $`C_f`$, $`b`$ and $`\theta `$ which are free parameters while the remaining variables are
$$\begin{array}{ccc}\hfill p& =& \frac{3\mathrm{\Lambda }\widehat{ϵ}\sqrt{\mathrm{\Lambda }(2m^23\mathrm{\Lambda })}}{m^26\mathrm{\Lambda }}\hfill \\ \hfill \lambda & =& \frac{8}{b^2}p\hfill \\ \hfill C_\varphi & =& \sqrt{\frac{m^26\mathrm{\Lambda }}{13p}}\hfill \\ \hfill C_A& =& \frac{mC_f}{C_h\left(p1\right)}\sqrt{\frac{13p}{m^26\mathrm{\Lambda }}}\hfill \\ \hfill C_{E(\pm )}& =& \frac{C_f}{\sqrt{2}}\sqrt{4\mathrm{\Lambda }p(m^2+6\mathrm{\Lambda })}\left(\frac{13p}{(m^26\mathrm{\Lambda })^3}\right)^{\frac{1}{4}}\hfill \end{array}$$
(3.25)
Here $`\theta =C_A_{\mathrm{}}`$ in (3.17). For $`\widehat{ϵ}=+1`$ and the particular case $`\mathrm{\Lambda }=m^2/6`$ corresponding to $`p=1/3\lambda =b^2/24`$ we have $`C_A=\sqrt{3}C_f/(2C_h)`$. The values of the remaining constants are well defined, $`C_\varphi =\sqrt{3}m`$ and $`C_{E(\pm )}=3C_fm^{5/2}/\sqrt{2}`$. For the values $`p=0`$ ($`\mathrm{\Lambda }=0`$) and $`p=1/2`$ ($`\mathrm{\Lambda }=m^2/2`$) we obtain $`C_E=0`$ and therefore the solutions presented here do not allow charged configurations for these particular limit values. In these cases $`C_\varphi m`$. For $`\widehat{ϵ}=1`$ the particular case $`\mathrm{\Lambda }=m^2/6`$ has no real solutions.
We have the bound in the cosmological constant
$$0<\mathrm{\Lambda }<\frac{m^2}{2}$$
(3.26)
such that $`p`$ is in the range
$$\{\begin{array}{ccccc}\hfill p& & ]0,\frac{1}{2}[\hfill & & \widehat{ϵ}=+1\hfill \\ \hfill p& & ]\mathrm{},0[]1,+\mathrm{}[\hfill & & \widehat{ϵ}=1\hfill \end{array}$$
(3.27)
where for both cases $`p=0`$ corresponds to $`\mathrm{\Lambda }=0`$ and for $`\widehat{ϵ}=+1`$ we have $`p=1/2`$ corresponding to $`x=\mathrm{\Lambda }/m^2=1/2`$ while for $`\widehat{ϵ}=1`$ we have $`p=1`$ corresponding to $`x=\mathrm{\Lambda }/m^2=1/2`$. For $`\widehat{ϵ}=+1`$ we have that $`p=1/3`$ corresponds to $`x=\mathrm{\Lambda }/m^2=1/6`$ while for $`\widehat{ϵ}=1`$ we have that $`lim_{x(1/6)^\pm }p=\mathrm{}`$. For $`\widehat{ϵ}=1`$, $`p]\mathrm{},0[`$ corresponds to $`x=\mathrm{\Lambda }/m^2]0,1/6[`$ and $`p]1,\mathrm{}[`$ corresponds to $`x=\mathrm{\Lambda }/m^2]1/6,1/2[`$.
### 3.2 Trivial Scalar Field Solutions: $`\varphi =0`$
It remains to analyse the case of $`\varphi =0`$. This case corresponds to not considering the scalar field at all and has been first addressed by Kogan , however in the original work a cosmological constant have not been considered (it has in but without solving the equations of motion), for this reason we also discuss it here.
Considering the above solutions for $`\gamma `$ (3.1), $`h`$ (3.2) and $`E`$ (3.3) the remaining three Einstein (A.27-A.29) reduce only to two independent equations
$`{\displaystyle \frac{m^2}{4}}+{\displaystyle \frac{\mathrm{\Lambda }}{2}}`$ $`=`$ $`\widehat{ϵ}\chi ^2`$ (3.28)
$`\alpha ^2+\alpha ^{}{\displaystyle \frac{3m^2}{4}}+{\displaystyle \frac{\mathrm{\Lambda }}{2}}`$ $`=`$ $`\widehat{ϵ}\chi ^2`$ (3.29)
while the scalar field equation (A.30) is already obeyed. Solving the first equation for $`\chi ^2`$ we get
$$\chi ^2=\widehat{ϵ}\left(\frac{m^2}{4}+\frac{\mathrm{\Lambda }}{2}\right)$$
(3.30)
and demanding the right hand side to be positive definite we obtain the constraint
$$\{\begin{array}{ccccc}\hfill \mathrm{\Lambda }& <& \frac{m^2}{2}\hfill & & \widehat{ϵ}=+1\hfill \\ \hfill \mathrm{\Lambda }& >& \frac{m^2}{2}\hfill & & \widehat{ϵ}=1\hfill \end{array}$$
(3.31)
As in the previous subsection in order to exist electric solutions the cosmological constant is constraint to be negative for $`\widehat{ϵ}=+1`$ and can be both negative in the range $`]m^2/2,0[`$ or positive for $`\widehat{ϵ}=1`$. We also note that from (3.30) the equality $`\mathrm{\Lambda }=m^2/2`$ holds that $`\chi =E=0`$, therefore not allowing electric configurations. For this reason we don’t consider the case $`\mathrm{\Lambda }=m^2/2`$.
From equation (3.29) and the definition of $`\alpha `$ (see (A.21) in appendix) we get the solution for $`f`$
$$f=c_f\mathrm{cosh}\left(\sqrt{k}(rr_0)\right)$$
(3.32)
where $`c_f`$ and $`r_0`$ are integration constants and
$$k=\frac{m^2}{2}\mathrm{\Lambda }.$$
(3.33)
From (3.2) we have that
$$h=c_h$$
(3.34)
and from (3.1) and the definition for $`\gamma `$ (A.21) we obtain the solution for $`A`$
$$A=\frac{mc_f\mathrm{sinh}\left(\sqrt{k}(rr_0)\right)}{c_h\sqrt{k}}+c_{A_0}$$
(3.35)
where $`c_{A_0}`$ is an integration constant that corresponds to the value of $`A`$ at $`r=r_0`$. Again we can set $`r_0=0`$ since it represents a shift in the radial coordinate.
For clarity we summarize and rewrite the solutions just obtained in the original frame
$$\begin{array}{ccc}\hfill h& =& C_h\hfill \\ \hfill f& =& C_f\mathrm{cosh}(Kr)\hfill \\ \hfill A& =& C_A\mathrm{sinh}(Kr)+\theta \hfill \\ \hfill E_{}& =& C_E\mathrm{cosh}(Kr)\hfill \\ \hfill A_0& =& \frac{C_E}{K}\mathrm{sinh}(Kr)\hfill \end{array}$$
(3.36)
where $`C_h`$ and $`C_f`$ are free constants and the remaining constants are defined as
$$\begin{array}{ccc}\hfill K& =& \sqrt{\frac{m^2}{2}\mathrm{\Lambda }}\hfill \\ \hfill C_A& =& \frac{mC_f}{C_hK}\hfill \\ \hfill C_{E(\pm )}& =& \pm \frac{C_f}{2}\sqrt{\left|\frac{m^2}{2}+\mathrm{\Lambda }\right|}\hfill \end{array}$$
(3.37)
The cosmological constant is constraint and accordingly $`K`$ is real for $`\widehat{ϵ}=+1`$
$$\widehat{ϵ}=+1:\{\begin{array}{ccc}\hfill \mathrm{\Lambda }& <& \frac{m^2}{2}\hfill \\ \hfill K& & ]m^2,+\mathrm{}[\hfill \end{array}$$
(3.38)
but can be both real and imaginary for $`\widehat{ϵ}=1`$
$$\widehat{ϵ}=1:\{\begin{array}{ccc}\hfill \mathrm{\Lambda }& & ]\frac{m^2}{2},0[\hfill \\ \hfill K& & ]0,m^2[\hfill \end{array}\mathrm{or}\{\begin{array}{ccc}\hfill \mathrm{\Lambda }& & \left]0.+\mathrm{}\right[\hfill \\ \hfill K& & ]0,+\mathrm{}[i\hfill \end{array}$$
(3.39)
where the last interval for $`K`$ is imaginary. In this last case we obtain periodic solutions in $`r`$ with period $`2\pi /|K|`$.
As a final remark we note that the contravariant electric density as defined in (2.19) is a constant
$$=\frac{C_EC_h}{C_f}$$
(3.40)
as expected from the solution for $`E`$ in the Cartan frame.
We already analyse the case for a null scalar field, but a constant scalar field is also an allowed trivial solution. For such solutions we obtain the same solutions up to the redefinition of the parameters
$$\stackrel{~}{\mathrm{\Lambda }}=\mathrm{\Lambda }e^{(ba)\varphi }\stackrel{~}{m}=me^{c\varphi }\stackrel{~}{\chi }=\chi e^{\frac{a}{2}c}$$
(3.41)
with $`\varphi =\mathrm{constant}`$.
## 4 Singularities, Geodesics and Horizons
### 4.1 Non-Trivial Solutions
The contraction of the Ricci tensor is
$$\begin{array}{ccc}\hfill R_{\mu \nu }R^{\mu \nu }& =& \frac{1}{4C_f^4r^4}[C_f^4(3+2p(8+p(17+2p(7+2p))))\hfill \\ & & 2(C_AC_fC_h(p1))^2(3+4p(p2))r+3(C_AC_h(p1))^4]\hfill \end{array}$$
(4.1)
which shows that there is a curvature singularity at $`r=0`$. The curvature is
$$R=\frac{h^3A^24f(hf^{\prime \prime }+f^{}h^{}+fh^{\prime \prime })}{2f^2h}=\frac{m^2p(34p)+6\mathrm{\Lambda }(4p^26p+1)}{2(m^26\mathrm{\Lambda })}\frac{1}{r^2}.$$
(4.2)
For $`\widehat{ϵ}=+1`$ we have always positive curvature while for $`\widehat{ϵ}=1`$ we can have both negative and positive curvatures
$$\begin{array}{ccccc}\hfill r^2R& & ]0,\frac{5}{8}[\hfill & \hfill \mathrm{for}& \widehat{ϵ}=+1\mathrm{and}x]0,\frac{1}{2}[\hfill \\ \hfill r^2R& & ]\mathrm{},0[\hfill & \hfill \mathrm{for}& \widehat{ϵ}=1\mathrm{and}x]0,\frac{9}{38}[/\left\{\frac{1}{6}\right\}\hfill \\ \hfill r^2R& =& 0\hfill & \hfill \mathrm{for}& \widehat{ϵ}=1\mathrm{and}x=\frac{9}{38}\hfill \\ \hfill r^2R& & ]0,\frac{9}{8}]\hfill & \hfill \mathrm{for}& \widehat{ϵ}=1\mathrm{and}x]\frac{9}{38},\frac{1}{2}[\hfill \end{array}$$
(4.3)
For the limiting cases $`\mathrm{\Lambda }0`$ (corresponding to $`x0`$) we have $`R0`$ for both $`\widehat{ϵ}=\pm 1`$ and for $`\mathrm{\Lambda }m^2/2`$ (corresponding to $`x1/2`$) we have $`R5/(8r^2)`$ for $`\widehat{ϵ}=+1`$ and $`R1/r^2`$ for $`\widehat{ϵ}=1`$. The Maximum value of the curvature for the case $`\widehat{ϵ}=1`$ is $`R=9/(8r^2)`$ corresponding to $`x=9/26`$. We note that as already explained in the last section $`x=0`$ and $`x=1/2`$ are not allowed solutions and, for $`\widehat{ϵ}=1`$, $`x=1/6`$ is neither an allowed solution.
For both cases $`\widehat{ϵ}=\pm 1`$ the curvature is asymptotically flat ($`lim_r\mathrm{}R=0`$), therefore our spaces are asymptotically flat.
In order to find if there is or not an horizon it is enough to consider a photon travelling in the radial direction. So we can solve equations (2.28) with $`L=0`$ and $`\kappa =0`$ obtaining
$$\begin{array}{ccc}\hfill t(r)& =& t_0\pm \frac{2}{|C_f|(2p3)}r^{3/2p}\hfill \\ \hfill \phi (r)& =& \phi _0\pm \frac{2\left((2p3)r^{p1}\theta \right)}{|C_f|(2p3)}r^{3/2p}\hfill \end{array}$$
(4.4)
For $`\widehat{ϵ}=+1`$ we have that $`p]0,1/2[`$, so these solutions are regular for all $`r`$ and we conclude that there is no horizon. From regularity at the singularity $`r=0`$ we are in the presence of a naked singularity, for an external observer the photon will hit the singularity in a finite time. For $`\widehat{ϵ}=1`$ we can have an horizon at $`r=0`$ as long as $`p>3/2`$ ($`p=3/2x=\mathrm{\Lambda }/m^2=9/26`$). This will happen for
$$x=\frac{\mathrm{\Lambda }}{m^2}]\frac{1}{6},\frac{9}{26}[.$$
(4.5)
Then in this range we will have a dressed singularity, for an external observer the infalling particle will take an infinite amount of time to reach the singularity. For all other values of $`p`$ we have a naked singularity. We note that from (4.4) for $`p=3/2`$ the geodesics are a fixed point on time and there are no horizons.
In order to understand the meaning of our singularity in terms of the angular variable let us now compute the angle deficit of our space, or equivalently the maximum value for the angular variable $`\phi `$. The metric reads
$$ds^2=r^{2p1}dt^2+dr^2+C_h^2r\left(d\phi +Adt\right)^2.$$
(4.6)
Let us remember from the discussion in section 2 that the $`2D`$ induced metric is $`h_{ij}=\mathrm{diag}(1,h^2)=\mathrm{diag}(1,C_h^2r)`$. Now let us make a transformation of coordinates $`r\stackrel{~}{r}`$ such that the measure of the induced metric is the usual one, i.e $`\sqrt{|h_{ij}|}=\stackrel{~}{r}^2`$. This accounts for a observer at rest in relation to space-time (hence rotating with space). The transformation of the radial coordinate is
$$r=\left(\frac{3}{4C_h}\right)^{\frac{2}{3}}\stackrel{~}{r}^{\frac{4}{3}}\{\begin{array}{ccc}\hfill f& =& \left(\frac{3\stackrel{~}{r}^2}{4C_h}\right)^{\frac{2}{3}\left(p\frac{1}{2}\right)}\hfill \\ \hfill h_{rr}& =& \left(\frac{4\stackrel{~}{r}}{3C_h^2}\right)^{\frac{2}{3}}\hfill \\ \hfill h_{\phi \phi }& =& \left(\frac{3C_h^2\stackrel{~}{r}^2}{4}\right)^{\frac{2}{3}}\hfill \end{array}$$
(4.7)
The maximum angle is computed as
$$\phi _{\mathrm{max}}=\frac{2\pi }{\sqrt{g}}\sqrt{\frac{h_{\phi \phi }}{h_{rr}}}=\frac{2\pi }{fh_{rr}}.$$
(4.8)
In order to obtain the background geometry we take the limit $`p0`$ (equivalent to $`\mathrm{\Lambda }0`$). We will discuss this limit properly in the next section when computing the mass, charge and angular momentum, for now let us just take it as granted, then the respective maximum angle is
$$\phi _{\mathrm{max}}=2\pi \frac{3|C_h|}{4}.$$
(4.9)
Imposing it to be as usal $`2\pi `$ we obtain the value for $`C_h`$
$$|C_h|=\frac{4}{3}.$$
(4.10)
So we have a rotating background without any angle deficit. For generic $`p`$ we obtain
$$\phi _{\mathrm{max}}=2\pi \left(\frac{3\stackrel{~}{r}}{4}\right)^{\frac{4}{3}p}$$
(4.11)
such that for $`\stackrel{~}{r}=4/3`$ we have $`\phi _{\mathrm{max}}=2\pi `$ always. In the limit $`\stackrel{~}{r}0`$ we obtain that for $`p>0`$, $`\phi _{\mathrm{max}}\mathrm{}`$ and for $`p<0`$, $`\phi _{\mathrm{max}}0`$. While in the limit $`\stackrel{~}{r}\mathrm{}`$ we obtain that for $`p>0`$, $`\phi _{\mathrm{max}}0`$ and for $`p<0`$, $`\phi _{\mathrm{max}}\mathrm{}`$. So we conclude that only for $`p<0`$ the singularity is a conical singularity (in the usual sense that we get an angular deficit), while for $`p>0`$ what we obtain as $`\stackrel{~}{r}0`$ is not a deficit, but instead a decompactification of the angular variable.
Then we have the following cases
$$\begin{array}{ccccc}\hfill \widehat{ϵ}=+1& x]0,\frac{1}{2}[\hfill & & p]0,\frac{1}{2}[:\hfill & \mathrm{decompactification}\mathrm{singularity}\hfill \\ \hfill \widehat{ϵ}=1& x]0,\frac{1}{6}[\hfill & & p]0,+\mathrm{}[:\hfill & \mathrm{decompactification}\mathrm{singularity}\hfill \\ \hfill \widehat{ϵ}=1& x]\frac{1}{6},\frac{1}{2}[\hfill & & p]\mathrm{},1[:\hfill & \mathrm{conical}\mathrm{singularity}\hfill \end{array}$$
(4.12)
### 4.2 Trivial Scalar Field Solutions
The contraction of the Ricci tensor is a constant
$$R_{\mu \nu }R^{\mu \nu }=\frac{K^4}{4C_f^4}\left(8C_f^48C_A^2C_f^2C_h^2+3C_A^4C_h^4\right)$$
(4.13)
which indicates that the space-time has no singularities. Specifically the curvature is
$$R=\frac{K^2}{2}\left(4\frac{C_A^2C_h^2}{C_f^2}\right)=\frac{m^2}{2}+2\mathrm{\Lambda }$$
(4.14)
and can have either positive or negative values. Taking in account the bounds for the cosmological constant (3.31) we obtain that
$$\begin{array}{ccccc}\hfill R& <& 0\hfill & \hfill \mathrm{for}& \widehat{ϵ}=+1\mathrm{and}\mathrm{\Lambda }<\frac{m^2}{2}\hfill \\ \hfill R& <& 0\hfill & \hfill \mathrm{for}& \widehat{ϵ}=1\mathrm{and}\mathrm{\Lambda }]\frac{m^2}{2},\frac{m^2}{4}[\hfill \\ \hfill R& =& 0\hfill & \hfill \mathrm{for}& \widehat{ϵ}=1\mathrm{and}\mathrm{\Lambda }=\frac{m^2}{4}\hfill \\ \hfill R& >& 0\hfill & \hfill \mathrm{for}& \widehat{ϵ}=1\mathrm{and}\mathrm{\Lambda }]\frac{m^2}{4},+\mathrm{}[.\hfill \end{array}$$
(4.15)
Therefore we conclude we are in the presence of an extended (non localized) configuration, there is no singularity, hence this solution cannot be considered as a classical particle. We recall that for $`\widehat{epsilon}=+1`$ and $`\mathrm{\Lambda }m^2/2`$ and for $`\widehat{epsilon}=1`$ and $`\mathrm{\Lambda }m^2/2`$ there are no allowed solutions.
## 5 Mass, Charge and Angular Momentum
In this section we compute the mass, charge and angular momentum.
### 5.1 Non-Trivial Scalar Field Solutions
As expected the Hamiltonian Constraint $`=0`$, Momentum Constraint $`^\phi =0`$ and Gauss Constraint $`𝒢=0`$ are obeyed, this is actually a way to check that our calculations are correct.
Using (2.23) we obtain that the Mass of the configuration is
$$\begin{array}{ccc}\hfill M& =& \left(2h^{}+4\lambda h\varphi \varphi ^{}\right)|_{r\delta _M}^r\mathrm{}=\frac{b^2+16\lambda }{b^2}C_h\frac{\mathrm{ln}(C_\varphi r)}{\sqrt{r}}|_{r\delta _M}^r\mathrm{}=\hfill \\ & =& 2C_hp\frac{\mathrm{ln}(C_\varphi \delta _M)}{\sqrt{\delta _M}}\hfill \end{array}$$
(5.1)
We introduced a cut-off $`\delta _M1`$ because this quantity has a infrared divergence as we compute the limit of $`\delta _M0`$.
The charge of this configuration is computed to be
$$Q_e=\frac{2C_hC_\varphi C_E}{C_f}$$
(5.2)
The constant $`C_f`$ can be set to unity by a proper redefinition of time $`tt/C_f`$ and the redefinitions of the remaining constants $`C_hC_h/C_f`$ and $`\theta \theta /C_f`$. So without any loss of generality we set $`C_f=1`$. However we must remember that $`C_E`$ as given in (3.25) has no defined sign and we must demand that the electric field has the correct sign when compared with the charge. From (5.2) we conclude that in order $`Q_e`$ and $`C_E`$ to have the same sign we are left only with the possibility of $`C_h<0`$, then $`C_h=4/3`$. Then we rewrite the charge as
$$Q_e=\pm \frac{2\sqrt{2}}{3}\frac{\sqrt{m^2p2\mathrm{\Lambda }(23p)}}{\left((m^26\mathrm{\Lambda })(13p)\right)^{\frac{1}{4}}}$$
(5.3)
where the $`\pm `$ accounts for positive and negative charge configurations. $`C_E`$ must account for this and the sign is set accordingly
$$C_E\mathrm{sign}(Q_e)$$
(5.4)
As we can see from (5.1) this choice of signal for $`C_h`$ affects the mass sign, the mass is positive or negative depending on the sign of $`p`$. We note that the logarithm in (5.1) is negative and therefore the mass is positive when $`p<0`$ and negative when $`p>0`$. For $`\widehat{ϵ}=+1`$ it is always negative, while for $`\widehat{ϵ}=1`$ it is negative for $`\mathrm{\Lambda }]0,m^2/6[`$ and positive for $`\mathrm{\Lambda }]m^2/6,m^2/2[`$
There is also one interesting point concerning the discrete symmetries time-inversion $`T`$ and $`P`$. Inverting time accounts for choosing $`C_f=1`$ such that $`tt`$. The visible direct effects of the transformation $`C_fC_f`$ for our solutions is to invert the sign of $`C_A`$ and $`C_E`$ (assuming we have fixed the $`\pm `$ of $`C_E`$, see (3.25)). Doing so we revert the sign of the charge definition as it depends explicitly on $`C_f`$ as well, see (5.2)), and although $`C_EC_E`$, our charge maintains its sign. Then we have two problems, first the charge and the electric field have now the wrong relative sign (we are considering $`C_h<0`$ fixed) and secondly the charge is not transforming properly under $`T`$ (see for instance equation (50) of , see also ). Therefore we are forced to transform $`C_hC_h`$ as well obtaining $`C_h>0`$. As a consequence $`C_A`$ does not actually changes sign (because the ratio $`C_h/C_f`$ does not change), this accounts fot $`T`$ violation due to the Chern-Simons term. Also we note that by choosing $`C_f=1`$ (or transforming $`C_fC_f`$ and $`C_hC_h`$) inverts the mass sign. This is actually expected, we recall the reader that classically a positron looks like an electron travelling backwards in time. As for parity $`P`$, will account for the transformation $`C_hC_h`$ which from the above discussion implies as well $`C_fC_f`$ and we obtain the same effects.
The Angular Momentum of this configuration is
$$J=\frac{2C_h^3C_A(p1)}{C_f}J_0=\frac{28m}{9}\sqrt{\frac{13p}{m^26\mathrm{\Lambda }}}J_0$$
(5.5)
where $`J_0`$ is the background angular momentum and will be computed later. the sign of $`J`$ does not depend in the particular configuration, but only on the relative sign between the Maxwell term ($`F^2`$) and the Chern-Simons term ($`AF`$) as explained on subsection 2.2. This means it will change if we consider the transformations $`mm`$ and vanishes for $`m=0`$ (as will be shown it does not vanishes in the limits $`m0^\pm `$, only for $`m=0`$). This is clearly also an effect of $`T`$ and $`P`$ violation which is expected when a Chern-Simons term is present.
So as we have just seen our solutions violate both $`T`$ and $`P`$ as expected when a Chern-Simons term is present. This is explicit on the fact that the signs of $`C_A`$ and $`J`$ only depend on the relative sign between the Maxwell and the Chern-Simons term.
We already computed the angle deficit in the last section such that for $`C_h=4/3`$ our background has the correct angular variable $`\phi [0,2\pi [`$. Here we still have to compute $`J_0`$, so we are properly explaining what are the limits of our solutions when we take the Chern-Simons coefficient to zero, $`m0`$. From the constraint interval we have that it corresponds to $`\mathrm{\Lambda }0`$ (equivalent to $`x=\mathrm{\Lambda }/m^20`$ and $`p0`$). We will analyse this limit from the definitions (3.25). In this limit we obtain from (3.25) that $`C_E0`$ therefore we have necessarily $`Q_e0`$, also we obtain $`C_\varphi 0`$ and $`C_A\mathrm{sign}(m)C_f/C_h`$. We note that for $`C_A`$ the limits on the right and left ($`m^\pm `$) are finite with opposite signs such that for $`x=p=0`$ we obtain $`C_A=0`$. Nevertheless the asymptotic limit are defined only from the left and from the right such that for the limiting cases $`C_A0`$. One obtains from (5.5) that $`J2\mathrm{sign}(m)C_h^2J_0`$. The first term corresponds to the background angular momentum, therefore we obtain
$$J_0=2\mathrm{sign}(m)C_h^2.$$
(5.6)
As already expected its sign depends on the relative sign between the Maxwell and the Chern-Simons term and accounts for parity violation. One obtains by a direct computation that $`M0`$ and also that the curvature vanishes everywhere, $`R0`$. Therefore as background for our configurations we obtain a stationary rotating flat space without any angle deficit as already studied in the last section. The background metric is
$$ds^2=\frac{1}{r}dt^2+dr^2+C_h^2r\left(d\phi +\left(\frac{\mathrm{sign}(m)}{C_h}\frac{1}{r}+\theta \right)dt\right)^2.$$
(5.7)
### 5.2 Trivial Scalar Field Solutions
We will now compute the charge, mass and angular momentum for the trivial solution (3.36) with $`\varphi =0`$.
The mass of the configuration is null, the charge is
$$Q_e=\frac{C_EC_h}{C_f}=\pm \frac{C_h}{2}\sqrt{\left|\frac{m^2}{2}+\mathrm{\Lambda }\right|}$$
(5.8)
and the angular momentum is
$$J=\frac{C_AC_h^3K}{C_f}J_0=mC_h^2J_0.$$
(5.9)
We note that again the sign of the angular momentum only depends on the relative sign between the Maxwell and Chern-Simons term.
We can solve (5.8) for $`C_h`$ obtaining
$$C_h=\frac{2Q_e}{\sqrt{\left|\frac{m^2}{2}+\mathrm{\Lambda }\right|}}$$
(5.10)
and
$$J=\frac{4mQ_e^2}{\left|\frac{m^2}{2}+\mathrm{\Lambda }\right|}J_0$$
(5.11)
Now the $`\pm `$ in $`C_E`$ must be chosen accordingly to the sign of the charge such that we obtain
$$C_E=\frac{\mathrm{sign}(Q_e)C_f}{2}\sqrt{\left|\frac{m^2}{2}+\mathrm{\Lambda }\right|}$$
(5.12)
Again we can redefine $`tt/C_f`$ that corresponds to set $`C_f=1`$.
By computing the limit $`m0`$ we obtain that $`C_A0`$, therefore both the charge and angular momentum vanish and we obtain the flat space
$$ds^2=C_fdt^2+dr^2+C_h(d\phi +\theta dt)^2.$$
(5.13)
Using the same procedure we obtain that the angular variable is in the range $`\phi [0,1/r^2[`$, so this space has some pathologies.
## 6 Summary and Discussion of Results
### 6.1 Summary of Non-Trivial Solutions
We will briefly resume the results obtained in this paper. Although we are repeating some of the equations of the article we think it is necessary in order to assemble and clarify all the results obtained.
We found a electric point particle that can constitute either a naked or dressed singularity, depending on the parameter choices. The results are presented in terms of $`x=\mathrm{\Lambda }/m^2`$, the cosmological constant to topological mass squared (Chern-Simons coefficient squared) ratio and the charge $`Q_e`$ of the configuration.
The metric, scalar field and gauge field solutions for such configuration are
$`\begin{array}{ccc}\hfill ds^2& =& \left({\displaystyle \frac{16}{9}}r\left(C_Ar^{p1}+\theta \right)^2r^{2p1}\right)dt^2+dr^2\hfill \\ & +& {\displaystyle \frac{16}{9}}rd\phi ^2+{\displaystyle \frac{16}{9}}r\left(C_Ar^{p1}+\theta \right)dtd\phi \hfill \\ \hfill \varphi & =& {\displaystyle \frac{2}{b}}\mathrm{ln}(|m|\sqrt{{\displaystyle \frac{16x}{13p}}}r)\hfill \\ \hfill A_0& =& {\displaystyle \frac{C_E}{p1}}r^{p1}\hfill \end{array}`$ (6.18)
where $`\theta `$ and $`b`$ are free parameters and all the remaining constants depend only on the cosmological constant to Chern-Simons square coefficient ratio $`x=\mathrm{\Lambda }/m^2`$
$`\begin{array}{ccc}\hfill p& =& {\displaystyle \frac{3x\widehat{ϵ}\sqrt{x(23x)}}{16x}}\hfill \\ \hfill C_A& =& {\displaystyle \frac{3\mathrm{sign}(m)}{4(1p)}}\sqrt{{\displaystyle \frac{(13p)}{(16x)}}}\hfill \\ \hfill C_E& =& {\displaystyle \frac{\mathrm{sign}(Q_e)\sqrt{p+2x(3p2)}}{\sqrt{2|m|}}}\left({\displaystyle \frac{13p}{(16x)^3}}\right)^{\frac{1}{4}}\hfill \end{array}`$ (6.22)
The Brans-Dicke coefficient is determined up to the free parameter $`b`$ as $`\lambda =8p/b^2`$ and the remaining scalar field exponential coefficients are fixed, $`a=0`$ and $`c=b/2`$.
The charge, angular momentum and mass are
$$\begin{array}{ccc}\hfill Q_e& =& \pm \frac{2\sqrt{2}}{3}\frac{\sqrt{p2x(23p)}}{\left((16x)(13p)\right)^{\frac{1}{4}}}\hfill \\ \hfill J& =& \frac{28\mathrm{sign}(m)}{9}\left(\sqrt{\frac{13p}{16x}}1\right)\hfill \\ \hfill M& =& \frac{8p}{3}\frac{\mathrm{ln}(C_\varphi \delta _M)}{\sqrt{\delta _M}}\frac{\mathrm{ln}\left(|m|\sqrt{\frac{16x}{13p}}\delta _M\right)}{\sqrt{\delta _M}}\hfill \end{array}$$
The mass is infrared divergent and we consider a cut-off proportional to the Planck Length, $`\delta _Ml_p=\sqrt{G}`$, being $`G`$ the Newton gravitational constant in natural units.
The curvature is
$$R=\frac{3\widehat{ϵ}\sqrt{x(23x)}+x(11+48x12\widehat{ϵ}\sqrt{x(23x)})}{2(16x)^2r^2}$$
(6.23)
and there is always a singularity at $`r=0`$ that we classify as decompactification or conical singularity depending if the range of $`\phi `$ goes to $`\mathrm{}`$ or $`0`$ (respectively) in the limit $`r0`$.
In the table below we present the possible ranges for $`x`$, $`p`$, $`\mathrm{\Lambda }`$, the sign of $`M`$ and the singularity classification. The $`\widehat{ϵ}`$ refers to the relative sign between the gauge sector and the gravitational sector.
$`\begin{array}{ccccccc}& & & & & & \\ & & & & & & \\ \widehat{ϵ}& x& p& \mathrm{\Lambda }& M& \mathrm{singularity}& \\ & & & & & & \\ & & & & & & \\ +1(ghosts)& ]0,1/2[& ]0,1/2[& ]0,m^2/2[& <0& decomp.& \\ 1(standard)& ]0,1/6[& ]\mathrm{},0[& ]0,m^2/6[& <0& decomp.& \\ & ]1/6,1/2[& ]1,+\mathrm{}[& ]m^2/6,m^2/2[& >0& conical& \end{array}`$ (6.30)
And to finalise we list the curvature sign and the existence or not of an horizon at $`r=0`$
$`\begin{array}{ccccc}& & & & \\ & & & & \\ \widehat{ϵ}& x& R& M& \mathrm{horizon}\\ & & & & \\ & & & & \\ +1(ghosts)& ]0,1/2[& <0& <0& no\\ 1(standard)& ]0,1/6[& <0& <0& no\\ & ]1/6,9/39[& <0& >0& yes\\ & ]9/39,9/26[& >0& >0& yes\\ & ]9/26,1/2[& >0& >0& no\end{array}`$ (6.39)
So we conclude that there is an horizon at $`r=0`$ only for standard fields and the $`x`$ range
$`x]{\displaystyle \frac{1}{6}},{\displaystyle \frac{9}{26}}[`$
such that we obtain a dressed singularity. We note that the mass of the solution are positive in this range. All remaining cases hold a naked singularity.
### 6.2 Summary of Trivial Solutions
We will summarize only the results for null scalar field ($`\varphi =0`$), i.e. solutions without the scalar field at all. This case have been addressed in without cosmological constant, we think is worthwhile to review these results with a non-null cosmological constant.
So, for $`\varphi =0`$ we found a electric extended configuration without singularities. The results are presented in terms of $`K=\sqrt{m^2/2\mathrm{\Lambda }}`$ and the charge $`Q_e`$ of the configuration. The metric and gauge field solutions for such configuration are
$`\begin{array}{cc}\hfill ds^2& =\left(\mathrm{cosh}^2(Kr)+C_h^2(C_A\mathrm{sinh}(Kr)+\theta )^2\right)dt^2+dr^2+C_h^2d\phi ^2\hfill \\ & +2C_h^2\left(C_A\mathrm{sinh}(Kr)+\theta \right)dtd\phi \hfill \\ \hfill A_0& ={\displaystyle \frac{C_E}{K}}\mathrm{sinh}(Kr)\hfill \end{array}`$ (6.43)
with
$`\begin{array}{ccc}\hfill K& =& \sqrt{{\displaystyle \frac{m^2}{2}}\mathrm{\Lambda }}\hfill \\ \hfill C_h& =& {\displaystyle \frac{2Q_e}{\sqrt{\left|\frac{m^2}{2}+\mathrm{\Lambda }\right|}}}\hfill \\ \hfill C_A& =& {\displaystyle \frac{m}{2Q_eK}}\sqrt{\left|{\displaystyle \frac{m^2}{2}}+\mathrm{\Lambda }\right|}\hfill \\ \hfill C_E& =& {\displaystyle \frac{\mathrm{sign}(Q_e)}{2}}\sqrt{\left|{\displaystyle \frac{m^2}{2}}+\mathrm{\Lambda }\right|}\hfill \end{array}`$ (6.48)
$`K`$ can be both real and imaginary. In the case of imaginary $`K`$ we obtain periodic solutions with period $`2\pi /|K|`$. We note that $`C_A`$ is multiplying by $`\mathrm{sinh}(Kr)`$ and correctly is also a pure imaginary such that $`g_{t\phi }`$ is real.
The mass of these configurations is null and the angular momentum is
$`J={\displaystyle \frac{4mQ_e^2}{\left|\frac{m^2}{2}+\mathrm{\Lambda }\right|}}.`$
There are no singularities and the curvature is constant
$`R=K^2+\mathrm{\Lambda }={\displaystyle \frac{m^2}{2}}+2\mathrm{\Lambda }.`$
In the next table we list the ranges for $`K`$, $`\mathrm{\Lambda }`$ and the sign of the curvature
$`\begin{array}{cccc}& & & \\ & & & \\ \widehat{ϵ}& \mathrm{\Lambda }& K& R\\ & & & \\ & & & \\ +1(ghosts)& ]m^2,+\mathrm{}[& ]\mathrm{},m^2/2[& <0\\ 1(standard)& ]0,m^2[& ]m^2/2,0[& <0\\ & ]0,+\mathrm{}[i& ]0,m^2/4]& 0\\ & ]9/26,1/2[& ]m^2/4,+\mathrm{}[& >0\end{array}`$ (6.56)
### 6.3 Discussion of Results
Given the Einstein Maxwell Chern-Simons theory coupled to a massless gravitational scalar field with action (2.1) discussed in section 2.1 we obtained the above classical solutions with electric charge only. We study both non-trivial and trivial solutions for the scalar field. For non-trivial solutions of the scalar field we obtain a rotating electric point particle that for the opposite sign between the gravitational and gauge sector and a certain range of the ratio $`\mathrm{\Lambda }/m^2`$ is dressed, while for trivial solutions of the scalar field we find an extend charge configuration that cannot be interpreted as a particle.
For non-trivial solutions it turns out that the solutions are highly constraint depending on the cosmological constant to Chern-Simons coefficient squared $`x=\mathrm{\Lambda }/m^2`$ which is constraint to the range $`x]0,1/2[`$. Further requiring that the background obtained (in the limit $`x0`$) to have no angular deficit we obtain only two free parameters, $`\theta `$ that accounts for the globally rotation of space and the $`\varphi `$ exponential coefficient $`b`$. Both of them are not relevant for any physical observables. We study both non-trivial and trivial solutions for the scalar field. Also we consider both the cases for the relative sign between the gauge sector and the gravitational sector $`\widehat{ϵ}=\pm 1`$. When they have the same sign ($`\widehat{ϵ}=+1`$) we have that the gauge fields are ghosts in the sense that contribute a negative amount of energy to the Hamiltonian, while in the case that they have opposite sign ($`\widehat{ϵ}=1`$) we have the standard case. Although the expressions for the solutions are expressed in the same way, the constants and consequently the physics change significantly. In particular the space-time curvature as well as the existence or non-existence of horizons will be sensitive to it.
For trivial solutions, the solutions are given in terms of $`K=\sqrt{m^2/2\mathrm{\Lambda }}`$ and the charge $`Q_e`$ and $`\theta `$ are free parameters. Although the cosmological constant is still bounded by the topological mass these bounds are not so restrictive. Again the relative sign $`\widehat{ϵ}=\pm 1`$ between the gravitational and gauge sector is relevant. In the limit $`m0`$ we obtain for $`\widehat{ϵ}=+1`$ that $`\mathrm{\Lambda }<0`$ while for $`\widehat{ϵ}=1`$ that $`\mathrm{\Lambda }>0`$. Our background is flat but with an angular deficit. In a similar way the curvature is sensitive to the relative sign $`\widehat{ϵ}`$.
The inclusion of the Chern-Simons topological term introduces very interesting features. Besides imposing the space to be rotating as explained in section 2.2 it imposes bounds on the cosmological constant trough the topological mass $`m`$. For the non-trivial solutions it constraints the allowed value for the cosmological constant to the interval $`\mathrm{\Lambda }]0,m^2/2[`$ such that the limit $`m0`$ corresponds also to $`\mathrm{\Lambda }0`$ (equivalent to $`x0`$ and $`p0`$) from the constraint $`0<\mathrm{\Lambda }<m^2/2`$ and we obtain in this limit a flat stationary background space-time. Then the cosmological constant is turn on and off by the Chern-Simons coefficient. It is very interesting that these facts emerges only as a consequence of the Chern-Simons term with out any ha-doc assumption. In this framework the existence of the cosmological constant can be interpreted as being due to the existence of the scalar field and the topological massive matter that constitute the electric point-particle. Therefore we can interpret that the charged matter deforms space-time such that the deformation is parameterized by the charge $`Q_e`$ and Brans-Dicke coefficient $`\lambda `$ and that the parameter $`x`$ is given as a function of $`Q_e`$ and $`\lambda `$. As expected this matter affects the curvature, either positively or negatively, depending on the sign of the gauge sector. For the trivial solution the cosmological constant bounds are not so restrictive but still exists a relation between topological mass and cosmological constant bounds, for $`\widehat{ϵ}=+1`$ we have $`\mathrm{\Lambda }<m^2/2`$ and for $`\widehat{ϵ}=1`$ we have $`\mathrm{\Lambda }>m^2/2`$.
As already mentioned, for non-trivial solutions, we have that the cosmological constant is always positive. However concerning the curvature we have different behaviours depending on the relative sign between the gauge and gravitational sector. For $`\widehat{ϵ}=+1`$ the curvature is always positive while for $`\widehat{ϵ}=1`$ the curvature is positive only for high values of $`x=\mathrm{\Lambda }/m^2]9/38,1/2[`$. To understand why let us contract the Einstein equations with the metric such that we obtain the relation
$`R=3e^{b\varphi }\mathrm{\Lambda }\lambda (\varphi ^2)+\widehat{ϵ}e^{c\varphi }E^2.`$
For solutions with $`\widehat{ϵ}=+1`$ the Brans-Dicke coefficient is always negative, hence all terms contribute positively to the curvature. For solutions with $`\widehat{ϵ}=+1`$ we have that the electric field contribution is always negative and that the Brans-Dicke coefficient is positive when $`x]0,1/6[`$ and negative when $`x]1/6,1/2[`$. Therefore we have the following cases, for $`x]0,1/6[`$ both the scalar field and electric field contribute negatively to the curvature while for $`x]1/6,1/2[`$ the scalar field contributes positively and the electric field contributes negatively to the curvature. We further note that from the expressions for the several constants (3.25) for the allowed solutions (3.24), we have that near $`x=1/6`$ the electric field contribution is predominant when compared with the scalar field contribution (that is negletable, $`C_\varphi 0`$). So only away from $`x=1/6`$ the scalar matter will become dominant over the charged matter and we have a positive curvature for $`x]9/38,1/2[`$. In this way we conclude that the scalar field is determinant in imposing the bounds on the cosmological constant (on the non-trivial solutions). Also it is the scalar field that allows for the existence of horizons. We concluded that there are horizons only for $`\widehat{ϵ}=1`$ in the range $`1/6<x<9/26`$ which corresponds to the greater positive values of $`p`$ ($`>3/2`$), remembering that the Brans-Dicke coefficient is proportional to $`\lambda p`$ this means that these values correspond to a region in which the scalar matter contributes positively to the curvature.
For the trivial solutions although the bound on the cosmological constant is not so restrictive the same behaviour concerning the curvature applies as can be seen directly in the expression for the curvature (4.14) that depends both in the cosmological constant and topological mass. For the trivial solutions we will have positive curvature only for $`\widehat{ϵ}=1`$ and $`\mathrm{\Lambda }>m^2/2`$.
The charges and angular momenta of the configurations are finite. The solutions are, for both non-trivial and trivial solutions of the scalar field, rotating spaces with angular momenta $`Jm`$ (or $`J\mathrm{sign}(m)`$), this accounts explicitly for the known parity $`P`$ and time-inversion $`T`$ violation due to the Chern-Simons term . Is explicit in the sense that the sign of the constant $`C_A`$ and of the angular momentum only depends on the relative sign between the Chern-Simons coefficient and the gravitational curvature term.
Concerning the mass of our configurations we concluded that its positiveness (or negativeness) is sensitive to the relative sign between the gravitational and gauge sector. However these results are not conclusive, although the charge and angular momentum are finite, the mass is infrared divergent, this is the main drawback of our solutions. The background is flat and therefore the reference mass (of the background) is null. Here we consider a cut-off of the order of the Planck Length. We believe that something is still missing in our theory, as already explained previously we are not considering a gravitational Chern-Simons. This correction to the Einstein action induces a correction to the configuration mass and would regularize it . For the extended trivial solutions the mass is null.
As a final remark we note that our solutions hold that $`a=0`$. Therefore the gravitational sector resembles an action with a dilatonic potential given by our cosmological constant term in (2.1), see for instance . We notice that by setting $`a=0`$ the field $`\varphi `$ is only minimally coupled to the $`2+1`$ metric and all the fields are expressed in terms of the scalar field (see the derivation of the solutions in section 3), therefore we would expect to obtain similar results by including more generic dilatonic potentials. An important point to stress here is that although our action is similar to the action of the work of Chan and Mann (CM) with an extra Chern-Simons term, it is not possible to obtain the solutions of those works in the limit $`m0`$. The main reason is that the although there the CM action is generic the authors only consider solutions for the particular case in which the scalar field can be interpreted as a dilaton. This means that the constants in our action (2.1) would be $`a=0`$, $`c=b=4`$ and $`\lambda =8`$ which is not the case since our constants are related as $`c=b/2`$ and $`\lambda `$ is dependent on several parameters. Therefore our massless scalar field cannot be interpreted as a dilaton. Secondly in our case we have no horizons away from $`r=0`$ and both our cosmological constant $`\mathrm{\Lambda }`$ and charge $`Q_e`$ vanishes in the limit $`m0`$ as already explained in detail, therefore we cannot possible obtain the solutions of Chan and Mann since their horizons are set uniquely by $`\mathrm{\Lambda }`$ and $`Q_e`$.
Interesting enough our gravitational field $`\varphi `$ can be related to the works of polarized cylindrical gravitational waves in $`3+1`$ gravity . For an explicit form of the effective $`2+1`$ dimensional action see for instance equation (1) of (see also ). In our case we further have a full gauge sector such that our classical solutions could constitute a possible electric charged background with cilindrical symmetry in $`3+1`$ dimensions (our solutions would correspond then to a electric charged string). Also similar actions have been considered in cosmological scenarios and in brane worlds .
After finishing this work the author realized that after we get our solutions redefining the radial coordinate accounts for changing the dilaton coupling (for $`a0`$) with the curvature $`R`$ and the Brans-Dicke parameter, however they will have generally different exponential factors, this does not invalidate the work presented here, simply we could yet consider a more generic action.
As an extension to this work the author intends to compute a pure magnetic solution using a similar action and procedure to this article. In order such configuration to exist it is necessary to consider an external electric charge distribution because as can be seen explicitly from Maxwell equations (2.2) or (A.25) for $`E=E_{}=0`$ we have that $`Bj^0`$. If we set $`j^0=0`$ the equations of motion hold that the magnetic field is null. This discussion is already put forward by Kogan (see conclusions of this reference). Another possible way out is to consider $`Ef=hAB`$ (such that $`E_{}=0`$, see (A.18) in appendix). In these cases the rotation will induce a electric field (see discussion in section 2.3). Also as other possible direction of research it would be interesting to consider extensions of this work that include gravitational Chern-Simons (as already explained we would expect to obtain finite mass) and dilatonic potentials.
Acknowledgments
This work was supported by Grant SFRH/BPD/5638/2001, SFRH/BPD/17683/2004 and POCTI/P-FIS-57547/2004. The author thanks both Ian Kogan for pointing out and Bayram Tekin for suggesting to further investigate gravitational solutions of Chern-Simons theories. The author thanks José Sande Lemos and Óscar Dias for invaluable discussions, also for explaining in detail part of their works as well recommending literature on the subject. The author thanks Nuno Reis and Carlos Herdeiro for several discussions and suggestions. The author thanks Paulo Vargas Moniz and Dmitri Gal’tsov for important remarks and suggestions.
## Appendix A Cartan Formalism
In this appendix we study the equations of motion in the Cartan Frame.
The Lagrangean 3-form corresponding to the action (2.1) is rewritten as
$$\begin{array}{cc}\hfill =& \{e^{a\varphi }[R1+2\lambda d\varphi d\varphi ]e^{b\varphi }\mathrm{\Lambda }1\hfill \\ & +\widehat{ϵ}e^{c\varphi }[FF+JA]+\widehat{ϵ}\frac{m}{2}AF\}\hfill \end{array}$$
(A.1)
with $`R`$ the metric curvature and $`F=dA`$ and where we define the Hodge dual as usual
$$(X)^{i_1\mathrm{}i_q}=(1)^D\frac{\sqrt{g}}{p!}ϵ^{i_1\mathrm{}i_qj_1\mathrm{}j_p}X_{j_1\mathrm{}j_p}$$
(A.2)
Introducing a triad $`\{e^0,e^1,e^2\}`$ such that
$$e^i=e_\alpha ^idx^\alpha g_{\alpha \beta }=\eta _{ij}e_\alpha ^ie_\beta ^j$$
(A.3)
where the Greek indices refer to the coordinates $`(x^0=t,x^1=r,x^2=\phi )`$ and the roman ones to the Cartan frame triad (meaning the flat space indices).
Varying the Lagrangean with respect to the Gauge field $`A`$, the coframe field $`e^i`$ and the dilaton $`\varphi `$ we obtain the equations of motion in the Cartan frame
$$\begin{array}{ccc}\hfill d(Fe^{c\varphi })J& =& \frac{m}{2}F\hfill \\ \hfill \left[e^{a\varphi }\left(G_{ij}+\mathrm{\Phi }_{ij}\right)e^{b\varphi }\eta _{ij}\mathrm{\Lambda }e_i2e^{c\varphi }T_{ij}\right]e^j& =& 0\hfill \\ \hfill e^{a\varphi }[(4a^2\lambda )dd\varphi +a(4a^22\lambda )d\varphi d\varphi ]& & \\ \hfill (b3a)e^{b\varphi }\mathrm{\Lambda }1& =& \widehat{ϵ}2(a+c)e^{c\varphi }FF\hfill \end{array}$$
(A.4)
respectively the Maxwell, Einstein and scalar field equations. We will specify the Einstein tensor $`G_{ij}`$, the Energy-Momentum tensor $`F_{ij}`$ and the scalar field tensor $`\mathrm{\Phi }_{ij}`$ for each metric parameterization used
$$\begin{array}{ccc}\hfill G_{ij}& =& R_{ij}\frac{1}{2}\eta _{ij}R\hfill \\ \hfill T_{ij}& =& \widehat{ϵ}\left(F_{ik}F_j^k\frac{1}{4}\eta _{ij}F_{kl}F^{kl}\right)\hfill \\ \hfill \mathrm{\Phi }_{ij}& =& a_i_j\varphi +a\eta _{ij}^2\varphi +(\lambda a^2)_i\varphi _j\varphi \left(\frac{\lambda }{2}a^2\right)\eta _{ij}_k\varphi ^k\varphi \hfill \end{array}$$
(A.5)
To proceed further one has to introduce a spin connection $`\omega _\alpha ^{ij}`$ and define the corresponding connection 1-form
$$\omega _j^i=\omega _{j\alpha }^idx^\alpha =\omega _{jk}^ie^k$$
(A.6)
Using the antisymmetric property (from definition)
$$\omega ^{ij}=\omega ^{ji}$$
(A.7)
and the Cartan Structure equation
$$de^i+\omega _j^ie^j=0$$
(A.8)
is enough to determine all the connection coefficients $`\omega _{jk}^i`$. In this work we are considering only radial symmetric configurations and metric parameterization such that $`e^1=dr`$ (note that a redefinition of $`r`$ introduces a non trivial metric component $`g_{11}`$) and $`e^0`$ and $`e^2`$ depend only on $`dt`$ and $`d\phi `$ (means that the metric has nonnull components $`g_{\alpha \alpha },g_{02}`$). In these particular cases we get the non vanishing connection coefficients
$$\begin{array}{c}\omega _{12}^0=\omega _{02}^1=\omega _{20}^1=\omega _{10}^2\omega _{21}^0=\omega _{01}^2\\ \omega _{10}^0=\omega _{00}^1\omega _{20}^0=\omega _{00}^2\\ \omega _{22}^0=\omega _{02}^2\omega _{22}^1=\omega _{12}^2\end{array}$$
(A.9)
plus the two equations
$$\begin{array}{c}de^0+\omega _{1}^0e^1+\omega _{2}^0e^2=0\\ de^2+\omega _{0}^2e^0+\omega _{1}^2e^1=0\end{array}$$
(A.10)
Also note that in this case the only Electric field component is $`E=F_{01}`$ ($`F_{02}=0`$ from Maxwell equations) and all the derivatives are with respect to $`r`$ only. Then it is now possible to define $`T_{ij}`$ and $`\mathrm{\Phi }_{ij}`$ for our parameterization:
$$\begin{array}{ccc}\hfill 2T_{00}& =& \widehat{ϵ}\left(B^2+E^2\right)\hfill \\ \hfill 2T_{11}& =& \widehat{ϵ}\left(B^2E^2\right)\hfill \\ \hfill 2T_{22}& =& \widehat{ϵ}\left(B^2+E^2\right)\hfill \\ \hfill 2T_{02}& =& 2\widehat{ϵ}BE\hfill \end{array}$$
(A.11)
the square of the Maxwell tensor is
$$F^2=2\widehat{ϵ}(B^2E^2)$$
(A.12)
and
$$\begin{array}{ccc}\hfill \mathrm{\Phi }_{00}& =& a\varphi ^{\prime \prime }+(\lambda /2a^2)(\varphi ^{})^2\hfill \\ \hfill \mathrm{\Phi }_{11}& =& \lambda /2(\varphi ^{})^2\hfill \\ \hfill \mathrm{\Phi }_{22}& =& a\varphi ^{\prime \prime }(\lambda /2a^2)(\varphi ^{})^2\hfill \end{array}$$
(A.13)
Note that the original electric field $`E_\alpha =F_{t\alpha }`$, magnetic field $`B_{}=F_{r\phi }`$ and external current $`J`$ are related to the Cartan frame ones $`E_i=F_{0i}`$, $`B=F_{12}`$ and $`j`$ either by using the triad $`e_\alpha ^i`$ or by the definition of the 2-forms $`F=F_{\alpha \beta }dx^\alpha dx^\beta =F_{ij}e^ie^j`$ and $`J=\sqrt{g}ϵ_{\mu \nu \rho }J^\mu dx^\nu dx^\rho =ϵ_{ijk}j^ie^je^k`$.
We use the metric parameterization such that the line element is given by
$$ds^2=f^2dt^2+dr^2+h^2(d\phi +Adt)^2$$
(A.14)
such that the usual components read
$$\begin{array}{ccc}\hfill g_{00}& =& f^2+h^2A^2\hfill \\ \hfill g_{11}& =& 1\hfill \\ \hfill g_{22}& =& h^2\hfill \\ \hfill g_{02}& =& h^2A\hfill \end{array}$$
(A.15)
The Cartan triad is then given by
$$\begin{array}{ccccccc}\hfill \mathrm{e}^0=d\theta ^0& =& fdt\hfill & & e_{0}^0=f\hfill & e_{1}^0=0\hfill & e_{2}^0=0\hfill \\ \hfill \mathrm{e}^1=d\theta ^1& =& dr\hfill & & e_{0}^1=0\hfill & e_{1}^1=1\hfill & e_{2}^1=0\hfill \\ \hfill \mathrm{e}^2=d\theta ^2& =& h(d\phi +Adt)\hfill & & e_{0}^2=hA\hfill & e_{1}^2=0\hfill & e_{2}^2=h\hfill \end{array}$$
(A.16)
such that the line element is now
$$ds^2=e^ie_i=\eta _{ij}d\theta ^id\theta ^j=(d\theta ^0)^2+(d\theta ^1)^2+(d\theta ^2)^2$$
(A.17)
with Minkowski metric $`\eta =\mathrm{diag}(1,1,1)`$.
The original Electric $`E_{}`$ and Magnetic fields $`B_{}`$ are given by
$$\begin{array}{ccc}\hfill E_{}& =& EfBhA\hfill \\ \hfill B_{}& =& Bh\hfill \end{array}$$
(A.18)
where $`E`$ and $`B`$ are the Cartan frame Electric and Magnetic fields.
The Cartan external currents $`j^i`$ are given by
$$\begin{array}{ccccc}\hfill j^0& =& fJ^t\hfill & =& \frac{e^{c\varphi }}{h}𝒥^t\hfill \\ \hfill j^2& =& h\left(J^\phi AJ^t\right)\hfill & =& \frac{e^{c\varphi }}{f}\left(𝒥^\phi A𝒥^t\right)\hfill \end{array}$$
(A.19)
where $`J^\mu `$ are the original external currents. For radial currents one has simply $`j^1=J^1=e^{c\varphi }𝒥^r/hf`$. We note that in terms of the physical $`𝒥^\mu `$ (measured by an external observer) we have $`J^\mu =e^{c\varphi }𝒥^\mu /hf`$ (see eq (2.20)).
From the form differentials
$$\begin{array}{ccc}\hfill de^0& =& \alpha e^0e^1\hfill \\ \hfill de^2& =& \beta e^1e^2\gamma e^0e^1\hfill \end{array}$$
(A.20)
we conclude that, except for the external currents, the Equations of motion, connections, curvature and so on depend only on the combinations
$$\alpha =\frac{f^{}}{f}\beta =\frac{h^{}}{h}\gamma =\frac{hA^{}}{f}$$
(A.21)
We list the non null connections in the Cartan frame
$$\begin{array}{c}\omega _{10}^0=\omega _{00}^1=\alpha \hfill \\ \omega _{12}^0=\omega _{21}^0=\omega _{02}^1=\omega _{20}^1=\omega _{01}^2=\omega _{10}^2=\gamma /2\hfill \\ \omega _{22}^1=\omega _{12}^2=\beta \hfill \end{array}$$
(A.22)
and the Einstein tensor components
$$\begin{array}{ccc}\hfill G_{00}& =& \beta ^2\gamma ^2/4\beta ^{}\hfill \\ \hfill G_{11}& =& \alpha \beta +\gamma ^2/4\hfill \\ \hfill G_{22}& =& \alpha ^23\gamma ^2/4+\alpha ^{}\hfill \\ \hfill G_{02}& =& \beta \gamma \gamma ^{}/2\hfill \end{array}$$
(A.23)
Then the Maxwell Equations are
$`B^{}+\alpha B+cB\varphi ^{}\gamma Ej^2`$ $`=`$ $`mEe^{c\varphi }`$ (A.24)
$`E^{}+\beta E+cE\varphi ^{}+j^0`$ $`=`$ $`mBe^{c\varphi }`$ (A.25)
The Einstein Equations are
$`e^{a\varphi }\left(\beta \gamma +{\displaystyle \frac{\gamma ^{}}{2}}\right)`$ $`=`$ $`2\widehat{ϵ}e^{c\varphi }EB`$ (A.26)
$`e^{a\varphi }\left[\beta ^2+{\displaystyle \frac{\gamma ^2}{4}}+\beta ^{}+a\varphi ^{\prime \prime }+\left(a^2{\displaystyle \frac{\lambda }{2}}\right)(\varphi ^{})^2\right]+{\displaystyle \frac{1}{2}}e^{b\varphi }\mathrm{\Lambda }`$ $`=`$ $`\widehat{ϵ}(B^2+E^2)e^{c\varphi }`$ (A.27)
$`e^{a\varphi }\left[\alpha ^2{\displaystyle \frac{3\gamma ^2}{4}}+\alpha ^{}+a\varphi ^{\prime \prime }+\left(a^2{\displaystyle \frac{\lambda }{2}}\right)(\varphi ^{})^2\right]+{\displaystyle \frac{1}{2}}e^{b\varphi }\mathrm{\Lambda }`$ $`=`$ $`\widehat{ϵ}(B^2+E^2)e^{c\varphi }`$ (A.28)
$`e^{a\varphi }\left[\alpha \beta +{\displaystyle \frac{\gamma ^2}{4}}+{\displaystyle \frac{\lambda }{2}}(\varphi ^{})^2\right]+{\displaystyle \frac{1}{2}}e^{b\varphi }\mathrm{\Lambda }`$ $`=`$ $`\widehat{ϵ}(B^2E^2)e^{c\varphi }`$ (A.29)
and the dilaton equation is
$$e^{a\varphi }\left[(4a^2\lambda )\varphi ^{\prime \prime }+a\left(4a^22\lambda \right)(\varphi ^{})^2\right]+(3ab)e^{b\varphi }\mathrm{\Lambda }=\widehat{ϵ}(a+c)(B^2E^2)e^{c\varphi }$$
(A.30) |
warning/0506/math0506067.html | ar5iv | text | # Small gaps between primes or almost primes
## 1. Introduction
In 1849, A. de Polignac (, p. 424) conjectured that every even number is the difference of two primes in infinitely many ways. More generally, we can let $`=\{h_1,h_2,\mathrm{},h_k\}`$ be a set of $`k`$ distinct integers. A major open question in number theory is to show that there are infinitely many positive integers $`n`$ such that $`n+h_1,n+h_2,\mathrm{},n+h_k`$ are all prime, provided that $``$ meets an obvious necessary condition that we call admissibility. For each prime $`p`$, let $`\nu _p()`$ be the number of distinct residue classes mod $`p`$ in $``$. We say that the set $``$ is admissible if $`\nu _p()<p`$ for all $`p`$.
Using heuristics from the circle method, Hardy and Littlewood realized the significance of the singular series $`𝔖()`$, defined as
(1.1)
$$𝔖()=\underset{p}{}\left(1\frac{\nu _p()}{p}\right)\left(1\frac{1}{p}\right)^k$$
for this problem. They made a conjecture about the asymptotic distribution of the numbers $`n`$ for which $`n+h_1,\mathrm{},n+h_k`$ are all prime, which we state here in the following form.
###### Conjecture 1.
Let $`\varpi (n)`$ denote function
(1.2)
$$\varpi (n)=\{\begin{array}{cc}\mathrm{log}n\hfill & \text{ if }n\text{ is prime,}\hfill \\ 0\hfill & \text{ otherwise.}\hfill \end{array}$$
As $`N`$ tends to infinity,
(1.3)
$$\underset{nN}{}\varpi (n+h_1)\varpi (n+h_2)\mathrm{}\varpi (n+h_k)=N(𝔖()+o(1)).$$
From the definition of $`𝔖()`$, we see that $`𝔖()0`$ if and only if $`\nu _p()<p`$ for all primes $`p`$; i.e., if and only if $``$ is admissible.
The set $`=\{0,2\}`$ is admissible, so the Hardy-Littlewood conjecture implies that
$$\underset{n\mathrm{}}{lim\; inf}(p_{n+1}p_n)=2,$$
where $`p_n`$ denotes the $`n^{\mathrm{th}}`$ prime. In an unpublished paper in the Partitio Numerorum series, Hardy and Littlewood proved that if the Generalized Riemann Hypothesis is true, then
$$\underset{n\mathrm{}}{lim\; inf}\left(\frac{p_{n+1}p_n}{\mathrm{log}p_n}\right)\frac{2}{3}.$$
In 1940, Erdős used Brun’s sieve to give the first unconditional proof of the inequality
$$\underset{n\mathrm{}}{lim\; inf}\left(\frac{p_{n+1}p_n}{\mathrm{log}p_n}\right)<1.$$
In 1965, Bombieri and Davenport proved unconditionally that
(1.4)
$$\underset{n\mathrm{}}{lim\; inf}\left(\frac{p_{n+1}p_n}{\mathrm{log}p_n}\right)0.4665\mathrm{}.$$
This result was one of the first applications of what is now known as the “Bombieri-Vinogradov Theorem,” which we state as follows.
###### Theorem (Bombieri-Vinogradov).
When $`(a,q)=1`$, let $`E(x;q,a)`$ be defined by the relation
(1.5)
$$\underset{\begin{array}{c}x<n2x\\ na(modq)\end{array}}{}\varpi (n)=\frac{x}{\varphi (q)}+E(x;q,a).$$
Furthermore, let
(1.6)
$$E(x,q)=\underset{a;(a,q)=1}{\mathrm{max}}|E(x,q,a)|,E^{}(N,q)=\underset{xN}{\mathrm{max}}E(x,q).$$
If $`A>0`$, then there exists $`B>0`$ such that if $`QN^{1/2}\mathrm{log}^BN`$, then
(1.7)
$$\underset{qQ}{}E^{}(N,q)_AN(\mathrm{log}N)^A.$$
This result was proved by Bombieri in 1965 (). At about the same time, A. I. Vinogradov () gave an independent proof of a slightly weaker result. There are numerous proofs of this result available in the literature; see, for example, and . We remark that in the usual definition of $`E(x;q,a)`$, one takes the sum in (1.5) to be over $`nx`$. However, the above definition is more convenient for our purposes.
The bound (1.4) was improved in several steps by Huxley to $`0.4394\mathrm{}`$. In 1988, Maier used his matrix method to improve the bound to $`0.2484\mathrm{}`$. Recently, the first, third and fourth authors proved a best possible result in this direction.
###### Theorem 1.
(Goldston, Pintz, and Yildirim )
$$\underset{n\mathrm{}}{lim\; inf}\left(\frac{p_{n+1}p_n}{\mathrm{log}p_n}\right)=0.$$
The proof of Theorem 1 uses, among other things, the Bombieri-Vinogradov Theorem. There are good reasons to believe that the bound in (1.7) holds for larger values of $`Q`$. More formally we have the following conjecture.
###### Hypothesis $`BV(\theta )`$.
Suppose $`1/2<\theta 1`$. If $`A>0,ϵ>0`$, then
(1.8)
$$\underset{qN^{\theta ϵ}}{}|E^{}(N;q,a)|_{A,ϵ}N(\mathrm{log}N)^A.$$
If Hypothesis $`BV(\theta )`$ is true, then we say that the sequence $`\varpi `$ has level of distribution $`\theta `$. Thus the Bombieri-Vinogradov Theorem shows that $`\varpi `$ has a level of distribution $`1/2.`$ The statement that $`\varpi `$ has a level of distribution $`1`$ is known as the “Elliott-Halberstam Conjecture” . Any level of distribution larger than $`1/2`$ will give the following strengthening of Theorem 1.
###### Theorem 2.
(Goldston, Pintz, and Yildirim ) If Hypothesis $`BV(\theta )`$ is true for some $`\theta >1/2`$, then
$$\underset{n\mathrm{}}{lim\; inf}(p_{n+1}p_n)<\mathrm{}.$$
If Hypothesis $`BV(\theta )`$ is true for some $`\theta `$ with $`4(8\sqrt{19})/15=0.97096\mathrm{}<\theta 1`$, then
$$\underset{n\mathrm{}}{lim\; inf}(p_{n+1}p_n)16.$$
Our first objective here is to give alternative proofs of Theorems 1 and 2. The primary difference in the proofs here and the proofs in comes from the use of Selberg diagonalization and a different choice of sieve coefficients; this will be discussed in more detail below. Our choice of coefficients allows us to give an elementary treatment of the main terms; we will discuss this further after the statement of Theorem 6 below.
Our second objective is to show that the results of can be strengthened if one replaces primes by numbers with a fixed number of prime factors. Let $`E_k`$ denote a number with numbers with exactly $`k`$ distinct prime factors. This contrasts with the usual definition of “almost-prime”, where $`P_k`$ is used to denote a number with at most $`k`$ distinct prime factors. Chen proved that there are infinitely many primes $`p`$ such that $`p+2`$ is a $`P_2.`$ While one expects that there are infinitely many primes $`p`$ such that $`p+2`$ is an $`E_2`$, this appears to be as difficult as the twin prime conjecture. However, we can prove that the limit infimum of gaps between $`E_2`$’s is bounded.
###### Theorem 3.
Let $`q_n`$ denote the $`n^{\mathrm{th}}`$ number that is a product of exactly two primes. Then
$$\underset{n\mathrm{}}{lim\; inf}\left(q_{n+1}q_n\right)26.$$
The above theorem uses an analogue of the Bombieri-Vinogradov theorem for the function $`\varpi \varpi `$, which is defined as
$$\varpi \varpi (n)=\underset{d|n}{}\varpi (d)\varpi (n/d).$$
Note that $`\varpi \varpi (n)=0`$ unless $`n`$ is a product of two primes or $`n`$ is a square of a prime.
When $`(a,r)=1`$, we have
$$\underset{\begin{array}{c}N<n2N\\ na(modr)\end{array}}{}\varpi \varpi (n)=\frac{1}{\varphi (r)}\underset{\chi (modr)}{}\overline{\chi }(a)\underset{N<n2N}{}\varpi \varpi (n)\chi (n),$$
and the expected value of this is
(1.9)
$$\frac{1}{\varphi (r)}\underset{N<n2N}{}\varpi \varpi (n)\chi _0(n),$$
where $`\chi _0`$ is the principal character mod $`r`$. A computation (see Lemma 7) shows that this quantity is asymptotically equal to
(1.10)
$$\frac{N}{\varphi (r)}\left(\mathrm{log}N+C_02\underset{p|r}{}\frac{\mathrm{log}p}{p}\right),$$
where $`C_0`$ is the absolute constant defined in (2.7).
Let $`E_2(N;r,a)`$ be defined by
$$\underset{\begin{array}{c}N<n2N\\ na(modr)\end{array}}{}\varpi \varpi (n)=\frac{N}{\varphi (r)}\left(\mathrm{log}N+C_02\underset{p|r}{}\frac{\mathrm{log}p}{p}\right)+E_2(N;q,a).$$
In parallel to the definitions of $`E(N,q)`$ and $`E^{}(N,q)`$, we define
$$E_2(N,r)=\underset{a,(a,r)=1}{\mathrm{max}}|E_2(N;r,a)|,E_2^{}(N,r)=\underset{xN}{\mathrm{max}}E_2(x,r).$$
###### Theorem (Bombieri-Vinogradov for $`\varpi \varpi `$).
For every $`A>0`$, there exists $`B>0`$ such that if $`QN^{1/2}\mathrm{log}^BN`$
$$\underset{rQ}{}|E_2^{}(N,r)|_AN(\mathrm{log}N)^A.$$
This is a special case of a result of Motohashi . Alternatively, one can easily modify the Vaughan’s Identity for the von Mangoldt function $`\mathrm{\Lambda }`$ to an identity for $`\mathrm{\Lambda }\mathrm{\Lambda }`$, and then use Vaughan’s approach (see or Chapter 28 of ) to the Bombieri-Vinogradov Theorem to prove the analogue for $`\mathrm{\Lambda }\mathrm{\Lambda }`$. It is then easy to modify this to a result for $`\varpi \varpi `$.
We also propose a natural analogue of Hypothesis $`BV(\theta )`$.
###### Hypothesis $`BV_2(\theta )`$.
Suppose $`1/2<\theta 1`$. If $`A>0,ϵ>0`$, then
(1.11)
$$\underset{qN^{\theta ϵ}}{}|E_2^{}(N;q)|_{A,ϵ}N(\mathrm{log}N)^A.$$
From this, we obtain the following conditional result.
###### Theorem 4.
If Hypotheses $`BV(\theta )`$ and $`BV_2(\theta )`$ are both true for some $`\theta `$ with $`(75\sqrt{473})/56=0.950918\mathrm{}<\theta 1`$, then
$$\underset{n\mathrm{}}{lim\; inf}\left(q_{n+1}q_n\right)6.$$
The basic construction for the proofs of Theorems 1 and 2 was inspired by work of Heath-Brown on almost prime-tuples of linear forms. Heath-Brown’s work was itself a generalization of Selberg’s proof that the polynomial $`n(n+2)`$ will infinitely often have at most five prime factors, and in such a way that one of $`n`$ and $`n+2`$ has at most two prime factors, while the other has at most three prime factors.
Define
(1.12)
$$P(n;)=\underset{h}{}(n+h),$$
The central idea is to relate the problem to sums of the form
(1.13)
$$\underset{N<n2N}{}\left(\underset{d|P(n;)}{}\lambda _d\right)^2$$
and of the form
(1.14)
$$\underset{N<n2N}{}\varpi (n)\left(\underset{d|P(n;)}{}\lambda _d\right)^2,$$
where one assumes that $`\lambda _d=0`$ for $`d>R`$, and $`R`$ is a parameter that is chosen to control the size of the error term. One also assumes that $`\lambda _d=0`$ when $`d`$ is not squarefree.
To illustrate the relevance of the sums (1.13) and (1.14), we discuss one simple application that is related to the second part of Theorem 2. Let $``$ be an admissible $`k`$-tuple, and consider the sum
(1.15)
$$𝒮:=\underset{N<n2N}{}\left\{\underset{h}{}\varpi (n+h)(\mathrm{log}3N)\right\}\left(\underset{d|P(n;)}{}\lambda _d\right)^2.$$
For a given $`n`$, the inner sum is negative unless there are at least two values $`h_i,h_j`$ such that $`n+h_i,n+h_j`$ are primes. From Theorems 5 and 6 below, one can deduce that if $`BV(\theta )`$ is true, if $`R=N^{\theta ϵ}`$ for $`ϵ>0`$, and if $`0\mathrm{}k`$, then
$$𝒮N𝔖()(\mathrm{log}R)^{k+2\mathrm{}}(\mathrm{log}N)m(k,\mathrm{},\theta ),$$
where
$$m(k,\mathrm{},\theta )=\left(\genfrac{}{}{0pt}{}{2\mathrm{}}{\mathrm{}}\right)\frac{1}{(k+2\mathrm{})!}\left\{\frac{k(2\mathrm{}+1)(\theta ϵ)}{(k+2\mathrm{}+1)(\mathrm{}+1)}1\right\}.$$
This last expression is positive, if for example, $`k=7`$, $`\mathrm{}=1`$, $`ϵ`$ is sufficiently small, and $`20/21<\theta 1`$. Consequently, if $`BV(1)`$ is true, then for any admissible $`7`$-tuple $``$, there are infinitely many $`n`$ and some $`h_i,h_j`$ such that $`n+h_i,n+h_j`$ are both prime. Now
$$=\{11,13,17,19,23,29,31\}$$
is an admissible $`7`$-tuple. $``$ is admissible because if $`p7`$, then none of the elements in $``$ are divisible by $`p`$, and if $`p>7`$, then there are not enough elements in $``$ to cover all of the residue classes mod $`p`$. Now any two elements of $``$ differ by at most $`20`$, so we conclude that if $`BV(1)`$ is true, then
$$\underset{n\mathrm{}}{lim\; inf}(p_{n+1}p_n)20.$$
To get the stronger bound of $`16`$ given in Theorem 2 needs an extra idea; this will be discussed in Section 7.
The success of the method depends upon making an appropriate choice for the $`\lambda _d`$, and this takes us into the realm of the Selberg upper bound sieve. It is a familiar fact from the theory of this sieve that
$$\underset{\begin{array}{c}N<n2N\\ d|P(n;)\end{array}}{}1=\frac{N}{f(d)}+r_d,$$
where $`f`$ is a multiplicative function and $`r_d`$ is a remainder term. (See the first part of Section 3 for the formal definition of $`f`$.) Accordingly, an appropriate transformation of the sum in (1.13) leads to consideration of the bilinear form
(1.16)
$$\underset{d,e}{}\frac{\lambda _d\lambda _e}{f([d,e])}.$$
The typical approach in the Selberg sieve is to choose the $`\lambda _d`$ to minimize the form in (1.16). To make this problem feasible, one needs to diagonalize this bilinear form. This can be done by making a change of variables
(1.17)
$$y_r=\mu (r)f_1(r)\underset{d}{}\frac{\lambda _{dr}}{f(dr)},$$
where $`f_1`$ is the multiplicative function defined by $`f_1=f\mu `$. (Note that the sum in (1.17) is finite because $`\lambda _d=0`$ for $`d>R`$.) The sum in (1.16) is then transformed into
$$\underset{r}{}\frac{y_r^2}{f_1(r)},$$
and the bilinear form is minimized by taking
(1.18)
$$y_r=\mu ^2(r)\frac{\lambda _1}{V},$$
where
$$V=\underset{r<R}{}\frac{\mu ^2(r)}{f_1(r)}.$$
The minimum of the form in (1.16) is then seen to be
$$\frac{\lambda _1^2}{V}.$$
One usually assumes that $`\lambda _1=1`$, but this is not an essential element of the Selberg sieve, and it is sometimes useful to assign some other nonzero value to $`\lambda _1`$.
The sum in (1.14) can be treated in a similar way. However, the corresponding function $`f`$ must be replaced by a slightly different function $`f^{}`$, which will be defined in Section 4. Therefore, the optimal choice of $`\lambda _d`$ is different from the optimal choice for the sum in (1.13). However, the basic structure of our approach requires that the same choice of $`\lambda _d`$ be used for both sums. We therefore face the problem of making a choice of $`\lambda _d`$ that works reasonably well for both problems. A similar choice was faced by Selberg and Heath-Brown in their earlier mentioned work, and they made this choice in different ways. Selberg made a choice of $`\lambda _d`$ that was optimal for one problem, and was able to successfully analyze the effect of this choice for the other problem. Heath-Brown chose
$$\lambda _d=\{\begin{array}{cc}\mu (d)\left(\frac{\mathrm{log}R/d}{\mathrm{log}R}\right)^{k+1}\hfill & \text{ if }d<R\text{,}\hfill \\ 0\hfill & \text{ otherwise;}\hfill \end{array}$$
$`k`$ being the number of linear forms under consideration. While this choice is not optimal for either problem, it is asymptotically optimal for both problems.
Inspired by Heath-Brown’s choice, Goldston, Pintz, and Yildirim chose
(1.19)
$$\lambda _{d,\mathrm{}}=\{\begin{array}{cc}\mu (d)\frac{(\mathrm{log}R/d)^{k+\mathrm{}}}{(k+\mathrm{})!}\hfill & \text{ if }d<R\text{,}\hfill \\ 0\hfill & \text{ otherwise.}\hfill \end{array}$$
Here, $`\mathrm{}`$ is a non-negative integer to be chosen in due course, with $`\mathrm{}k`$. With the exponent $`k+\mathrm{}`$, one is effectively using a $`k+\mathrm{}`$-dimensional sieve on a $`k`$-dimensional sieve problem. In an upper bound sieve, it is optimal to take the dimension of the sieve to be the same as the dimension of the problem. In the problems considered here, however, it is not the upper bound but the ratio of the quantities in (1.13) and (1.14) that is relevant. The presence of the parameter $`\mathrm{}`$ is essential for the success of their method.
In the current exposition, we make a choice that is a hybrid of the above and Selberg’s original approach. Our choice is most easily described in terms of $`y_r`$. We choose
(1.20)
$$y_{r,\mathrm{}}=y_{r,\mathrm{}}()=\{\begin{array}{cc}\frac{\mu ^2(r)𝔖()(\mathrm{log}R/r)^{\mathrm{}}}{\mathrm{}!}\hfill & \text{ if }r<R\text{,}\hfill \\ 0\hfill & \text{ otherwise.}\hfill \end{array}$$
As motivation for this choice, we note that $`y_{r,0}`$ is the optimal choice given in (1.18) with $`\lambda _1=V𝔖()`$. Moreover, one can show that
$$\mu (r)f_1(r)\underset{d<R/r}{}\frac{\mu (dr)}{f(dr)}\frac{\mathrm{log}^{k+\mathrm{}}(R/rd)}{(k+\mathrm{})!}\frac{𝔖()(\mathrm{log}R/r)^{\mathrm{}}}{\mathrm{}!}$$
when $`r`$ is not too close to $`R`$. In other words, the choice of $`\lambda _{d,\mathrm{}}`$ in (1.19) gives a value of $`y_r`$ that is asymptotic to the expression in (1.20).
One can use (1.17) and Möbius inversion to deduce that
(1.21)
$$\frac{\lambda _{d,\mathrm{}}}{f(d)}=\mu (d)\underset{r}{}\frac{y_{dr,\mathrm{}}}{f_1(rd)},$$
and so, when the choice of $`y_{r,\mathrm{}}`$ of (1.20) is specified, one obtains
(1.22)
$$\lambda _{d,\mathrm{}}=\mu (d)\frac{f(d)}{f_1(d)}\frac{𝔖()}{\mathrm{}!}\underset{\begin{array}{c}r<R/d\\ (r,d)=1\end{array}}{}\frac{\mu ^2(r)}{f_1(r)}(\mathrm{log}R/rd)^{\mathrm{}}$$
when $`d<R`$. With this choice of $`\lambda _{d,\mathrm{}}`$, we set
(1.23)
$$\mathrm{\Lambda }_R(n;,\mathrm{})=\underset{d|P(n;)}{}\lambda _{d,\mathrm{}}.$$
As we shall see, this choice $`\lambda _{d,\mathrm{}}`$ allows us to give elementary estimates for the main terms in (1.13) and (1.14).
We also define
(1.24)
$$\beta ()=\underset{p}{}\frac{(k\nu _p())\mathrm{log}p}{p}.$$
This sum is finite because $`\nu _p=k`$ for sufficiently large $`p`$.
Theorems 1 through 4 will be derived fairly easily from the following results.
###### Theorem 5.
Suppose that $`=\{h_1,\mathrm{},h_k\}`$ is an admissible set, and that $`0\mathrm{}_1,\mathrm{}_2k`$. If $`RN^{1/2ϵ}`$ then
(1.25) $`{\displaystyle \underset{N<n2N}{}}`$ $`\mathrm{\Lambda }_R(n;,\mathrm{}_1)\mathrm{\Lambda }_R(n;,\mathrm{}_2)=`$
$`\left({\displaystyle \genfrac{}{}{0pt}{}{\mathrm{}_1+\mathrm{}_2}{\mathrm{}_1}}\right)𝔖()N{\displaystyle \frac{(\mathrm{log}R)^{k+\mathrm{}_1+\mathrm{}_2}}{(k+\mathrm{}_1+\mathrm{}_2)!}}\left\{1+O(\beta ()𝔖()/\mathrm{log}R)\right\}.`$
The implied constant depends at most on $`k`$.
###### Theorem 6.
Suppose that $`=\{h_1,\mathrm{},h_k\}.`$ Suppose further that Hypothesis $`BV(\theta )`$ is true and $`RN^{(\theta ϵ)/2}`$. If $`h_0`$, $``$ is admissible, and $`0\mathrm{}_1,\mathrm{}_2k`$, then
(1.26) $`{\displaystyle \underset{N<n2N}{}}`$ $`\varpi (n+h_0)\mathrm{\Lambda }_R(n;,\mathrm{}_1)\mathrm{\Lambda }_R(n;,\mathrm{}_2)=`$
$`\left({\displaystyle \genfrac{}{}{0pt}{}{\mathrm{}_1+\mathrm{}_2+2}{\mathrm{}_1+1}}\right)N𝔖(){\displaystyle \frac{(\mathrm{log}R)^{k+\mathrm{}_1+\mathrm{}_2+1}}{(k+\mathrm{}_1+\mathrm{}_2+1)!}}\left\{1+O(\beta ()𝔖()/\mathrm{log}R)\right\}.`$
If $`h_0`$, $`^0=\{h_0\}`$ is admissible, and $`1\mathrm{}_1,\mathrm{}_2k`$ then
(1.27) $`{\displaystyle \underset{N<n2N}{}}`$ $`\varpi (n+h_0)\mathrm{\Lambda }_R(n;,\mathrm{}_1)\mathrm{\Lambda }_R(n;,\mathrm{}_2)=`$
$`\left({\displaystyle \genfrac{}{}{0pt}{}{\mathrm{}_1+\mathrm{}_2}{\mathrm{}_1}}\right)N𝔖(^0){\displaystyle \frac{(\mathrm{log}R)^{k+\mathrm{}_1+\mathrm{}_2}}{(k+\mathrm{}_1+\mathrm{}_2)!}}\{1+O(\beta (^0)𝔖(^0)/\mathrm{log}R\}.`$
The implied constants depend at most on $`k`$.
With a bit more work, we could allow $`\mathrm{}_1`$ or $`\mathrm{}_2`$ to be $`0`$ in (1.27). However, we omit this because the only place we use this result is in the proof of Theorem 1, where we will have $`\mathrm{}_1=\mathrm{}_2>0`$.
Analogues of Theorems 5 and 6 are given in for $`\lambda _{d,\mathrm{}}`$ given by (1.19). The corresponding main terms in are evaluated with the help of contour integrals in two variables and zero-free regions for the Riemann-zeta function. On the other hand, with the choice of $`\lambda _{d,\mathrm{}}`$ given in (1.22), we are able to give an elementary treatment of the main terms in Theorems 5 and 6.
###### Theorem 7.
Suppose that $`=\{h_1,\mathrm{},h_k\}`$ is an admissible set, and that $`0\mathrm{}_1,\mathrm{}_2k`$. Suppose that Hypotheses $`BV(\theta )`$ and $`BV_2(\theta )`$ are both satisfied, and $`RN^{(\theta ϵ)/2}`$. If $`h_0`$, then
$`{\displaystyle \underset{N<n2N}{}}`$ $`\varpi \varpi (n+h_0)\mathrm{\Lambda }_R(n;,\mathrm{}_1)\mathrm{\Lambda }_R(n;,\mathrm{}_2)=`$
$`\{\left({\displaystyle \genfrac{}{}{0pt}{}{\mathrm{}_1+\mathrm{}_2+2}{\mathrm{}_1+1}}\right)(N\mathrm{log}N)𝔖(){\displaystyle \frac{(\mathrm{log}R)^{k+\mathrm{}_1+\mathrm{}_2+1}}{(k+\mathrm{}_1+\mathrm{}_2+1)!}}`$
$`+2T(k,\mathrm{}_1,\mathrm{}_2)N𝔖(){\displaystyle \frac{(\mathrm{log}R)^{k+\mathrm{}_1+\mathrm{}_2+2}}{(k+\mathrm{}_1+\mathrm{}_2+2)!}}\left\}\right\{1+O(\beta ()𝔖()/\mathrm{log}R)\},`$
where
$$T(k,\mathrm{}_1,\mathrm{}_2)=\left(\genfrac{}{}{0pt}{}{\mathrm{}_1+\mathrm{}_2+3}{\mathrm{}_2+1}\right)\left(\genfrac{}{}{0pt}{}{\mathrm{}_1+\mathrm{}_2+3}{\mathrm{}_1+1}\right)+\left(\genfrac{}{}{0pt}{}{\mathrm{}_1+\mathrm{}_2+2}{\mathrm{}_1+1}\right).$$
The implied constant depends at most on $`k`$.
The reader will note that the sums considered here are more general than the sums in (1.13) and (1.14)–the latter correspond to the case $`\mathrm{}_1=\mathrm{}_2=\mathrm{}`$. We will see in Section 7 that this extra flexibility is useful in applications.
We also remark that the proof of Theorem 1 requires averaging over a set of $``$, where the elements of $``$ can be as large as $`\mathrm{log}R`$. Accordingly, we shall take some extra effort to make our estimates uniform in $`h`$ under the assumption that $`h\mathrm{log}N`$. For our results, it is not necessary to make the estimates in Theorems 5 through 7 uniform in $`k`$.
The implied constants in the error terms of Theorems 6 and 7 are ineffective due to the use of the Bombieri-Vinogradov Theorem, which uses the Siegel-Walfisz Theorem. However, the constants can be made effective by using the procedure of Section 13 of . This procedure deletes the greatest prime factor of the eventually existing exceptional modulus from the sieve process.
The paper gives an unconditional proof of a quantitative version of Theorem 1; that
(1.28)
$$\underset{n\mathrm{}}{lim\; inf}\frac{(p_{n+1}p_n)}{\mathrm{log}p_n(\mathrm{log}\mathrm{log}p_n)^1\mathrm{log}\mathrm{log}\mathrm{log}\mathrm{log}p_n}<\mathrm{},$$
and this result requires that the estimates in Theorems 5 and 6 be uniform in $`k`$. In a forthcoming paper, Goldston, Pintz, and Yildirim will improve (1.28) to
(1.29)
$$\underset{n\mathrm{}}{lim\; inf}\frac{(p_{n+1}p_n)}{(\mathrm{log}p_n)^{1/2}(\mathrm{log}\mathrm{log}p_n)^2}<\mathrm{}.$$
The function $`\varpi \varpi `$ used in Theorem 7 is convenient for calculations, but it is not optimal for applications. In a future paper we will show that by using other functions supported on $`E_2`$’s, the bound in Theorem 3 can be improved to 8 and the allowable range for $`\theta `$ in Theorem 4 can be improved to $`0.51<\theta 1`$. We will also show that there is a constant $`C`$ such that for any positive integer $`r`$,
$$\underset{n\mathrm{}}{lim\; inf}(q_{n+r}q_n)Cre^r.$$
Notation: The letters $`R,N`$ denote real variables tending to infinity. The letter $`p`$ is always used to denote a prime. The letters $`d,e,r`$ are usually squarefree numbers; the letters $`m,n`$ are usually positive integers. The notation $`\omega (n)`$ is used to denote the number of distinct prime factors of $`n`$. We use $`\rho `$ to denote the function
$$\rho (r)=1+\underset{p|r}{}\frac{\mathrm{log}p}{p}.$$
The letters $`S,,U,`$ and $`V`$, with or without subscripts, are often used to denote sums. The meanings of these symbols are local to sections; e.g., the meaning of $`S_1`$ in Section 6 is different from the meaning of $`S_1`$ in Section 7.
We use $`_{}^{\mathrm{}}{}_{}{}^{}`$ to denote a summation over squarefree integers. In general, the constants implied by “$`O`$” and “$``$” will depend on $`k`$. Any other dependencies will be explicitly noted. As noted before, $`k`$ is the size of $``$; we always assume that $`k2`$. The parameter $`\mathrm{}`$, with or without subscript, is an integer with $`0\mathrm{}k`$.
## 2. Preliminary Lemmas
The following two lemmas are classical estimates that have proved useful for handling remainder terms that arise in the Selberg sieve. The results can be found in Halberstam and Richert’s book (, Lemmas 3.4 and 3.5). We reproduce the proofs here since they are quite short.
###### Lemma 1.
For any natural number $`h`$ and for $`x1`$,
$`{\displaystyle \stackrel{\mathrm{}}{}_{dx}}{\displaystyle \frac{h^{\omega (d)}}{d}}`$ $`(\mathrm{log}x+1)^h,`$
$`{\displaystyle \stackrel{\mathrm{}}{}_{dx}}h^{\omega (d)}`$ $`x(\mathrm{log}x+1)^h.`$
###### Proof.
For the first inequality, we note that the sum on the left is
$$\underset{d_1\mathrm{}d_hx}{}\frac{\mu ^2(d_1\mathrm{}d_h)}{d_1\mathrm{}d_h}\left(\underset{nx}{}\frac{1}{n}\right)^h(\mathrm{log}x+1)^h.$$
For the second inequality, we note that the left-hand side is at most
$$x\stackrel{\mathrm{}}{}_{dx}\frac{h^{\omega (d)}}{d},$$
and we appeal to the first inequality. ∎
###### Lemma 2.
Assume Hypothesis $`BV(\theta )`$, and let $`h`$ be a positive integer. Given any positive constant $`U`$ and any $`ϵ>0`$, then
$$\stackrel{\mathrm{}}{}_{d<N^{\theta ϵ}}h^{\omega (d)}E^{}(N,d)_{U,h,ϵ}N(\mathrm{log}N)^U.$$
Similarly, if Hypothesis $`BV_2(\theta )`$ is assumed, then
$$\stackrel{\mathrm{}}{}_{d<N^{\theta ϵ}}h^{\omega (d)}E_2^{}(N,d)_{U,h,ϵ}N(\mathrm{log}N)^U.$$
###### Proof.
We begin by noting the trivial estimate $`E^{}(N,d)N(\mathrm{log}N)/d`$. By Cauchy’s inequality
$`{\displaystyle \stackrel{\mathrm{}}{}_{d<N^{\theta ϵ}}}h^{\omega (d)}E^{}(N,d)`$ $`\left(N\mathrm{log}N{\displaystyle \stackrel{\mathrm{}}{}_{d<N^{\theta ϵ}}}{\displaystyle \frac{h^{2\omega (d)}}{d}}\right)^{1/2}\left({\displaystyle \stackrel{\mathrm{}}{}_{d<N^{\theta ϵ}}}E^{}(N,d)\right)^{1/2}`$
$`_{h,ϵ,A}N(\mathrm{log}N)^{(h^2A+1)/2}.`$
We have used Lemma 1 and Hypothesis $`BV(\theta )`$ in the last line. The first result follows by taking $`A=h^2+1+2U`$. The second result is proved similarly; one uses the trivial bound $`E_2^{}(N,d)N(\mathrm{log}N)^2/d`$. ∎
###### Lemma 3.
If $`a,b`$ are positive real numbers, both at least $`1`$, then
$$_1^x(\mathrm{log}x/u)^{a1}(\mathrm{log}u)^{b1}\frac{du}{u}=(\mathrm{log}x)^{a+b1}\frac{\mathrm{\Gamma }(a)\mathrm{\Gamma }(b)}{\mathrm{\Gamma }(a+b)}.$$
###### Proof.
Upon making the change of variables $`u=x^v`$, the left-hand side becomes
$$(\mathrm{log}x)^{a+b1}_0^1(1v)^{a1}v^{b1}𝑑v.$$
The result follows by the standard formula for the beta-integral. ∎
Our next lemma is another standard result in the theory of sieves.
###### Lemma 4.
Suppose that $`\gamma `$ is a multiplicative function, and suppose that there positive real numbers $`\kappa ,A_1,A_2,L`$ such that
(2.1)
$$0\frac{\gamma (p)}{p}1\frac{1}{A_1},$$
and
(2.2)
$$L\underset{wp<z}{}\frac{\gamma (p)\mathrm{log}p}{p}\kappa \mathrm{log}\frac{z}{w}A_2$$
if $`2wz`$. Let $`g`$ be the multiplicative function defined by
(2.3)
$$g(d)=\underset{p|d}{}\frac{\gamma (p)}{p\gamma (p)}.$$
Then
$$\stackrel{\mathrm{}}{}_{d<z}g(d)=c_\gamma \frac{(\mathrm{log}z)^\kappa }{\mathrm{\Gamma }(\kappa +1)}\left\{1+O_{A_1,A_2,\kappa }\left(\frac{L}{\mathrm{log}z}\right)\right\},$$
where
$$c_\gamma =\underset{p}{}\left(1\frac{\gamma (p)}{p}\right)^1\left(1\frac{1}{p}\right)^\kappa .$$
This is a combination of Lemmas 5.3 and 5.4 of Halberstam and Richert’s book . In , the hypothesis (2.1) is denoted $`(\mathrm{\Omega }_1)`$, and hypothesis (2.2) is denoted $`(\mathrm{\Omega }_2(\kappa ,L))`$. As indicated above, the constant implied by “$`O`$” may depend on $`A_1,A_2,\kappa `$, but it is independent of $`L`$. This will be important in our applications.
###### Lemma 5.
Suppose that $`\gamma `$ and $`g`$ satisfy the same hypotheses as in the previous lemma. If $`a`$ is a non-negative integer, then
$$\stackrel{\mathrm{}}{}_{r<R}g(r)(\mathrm{log}R/r)^a=c_\gamma \frac{\mathrm{\Gamma }(a+1)}{\mathrm{\Gamma }(\kappa +a+1)}(\mathrm{log}R)^{\kappa +a}+O_{A_1,A_2,\kappa ,a}\left(L(\mathrm{log}R)^{\kappa +a1}\right).$$
###### Proof.
When $`a=0`$, this is Lemma 4. If $`a>0`$, then
$`{\displaystyle \stackrel{\mathrm{}}{}_{r<R}}g(r)(\mathrm{log}R/r)^a`$ $`=a{\displaystyle \stackrel{\mathrm{}}{}_{r<R}}g(r){\displaystyle _r^R}(\mathrm{log}R/z)^{a1}{\displaystyle \frac{dz}{z}}`$
$`={\displaystyle _1^R}{\displaystyle \frac{a(\mathrm{log}R/z)^{a1}}{z}}{\displaystyle \stackrel{\mathrm{}}{}_{r<z}}g(r)dz.`$
Using Lemma 4, we see that the above is
$`{\displaystyle _1^R}`$ $`{\displaystyle \frac{a(\mathrm{log}R/z)^{a1}}{z}}\left\{{\displaystyle \frac{c_\gamma (\mathrm{log}z)^\kappa }{\mathrm{\Gamma }(\kappa +1)}}+O(L(\mathrm{log}z)^{\kappa 1})\right\}dz`$
$`=`$ $`{\displaystyle \frac{ac_\gamma }{\mathrm{\Gamma }(\kappa +1)}}{\displaystyle _1^R}(\mathrm{log}R/z)^{a1}(\mathrm{log}z)^\kappa {\displaystyle \frac{dz}{z}}+O\left(aL{\displaystyle _1^R}(\mathrm{log}R/z)^{a1}(\mathrm{log}z)^{\kappa 1}{\displaystyle \frac{dz}{z}}\right).`$
The desired result follows from using Lemma 3. ∎
###### Lemma 6.
If $``$ is admissible and $`|h_i|h`$ for all $`h_i`$, then
(2.4)
$$1\beta ()\mathrm{log}\mathrm{log}10h$$
and there is a constant $`b_k`$ (depending only on $`k`$) such that
(2.5)
$$𝔖()(\mathrm{log}\mathrm{log}10h)^{b_k}.$$
###### Proof.
Without loss of generality, we may assume that $`h100`$; this will simplify the writing of logarithms. We note that $`\nu _p<k`$ if and only if $`p|\mathrm{\Delta }()`$, where
(2.6)
$$\mathrm{\Delta }=\mathrm{\Delta }():=\underset{1i<jk}{}|h_ih_j|.$$
Therefore
$$\beta ()=\underset{p|\mathrm{\Delta }}{}(k\nu _p)\frac{\mathrm{log}p}{p},$$
where we have written $`\nu _p`$ as an abbreviation for $`\nu _p()`$. We may assume without loss of generality that $`\mathrm{\Delta }100`$.
Now $`\nu _2=1`$ whenever $``$ is admissible, so we see that $`\beta ()\mathrm{log}2/2`$. In the opposite direction, we have
$`\beta ()`$ $`{\displaystyle \underset{p\mathrm{log}\mathrm{\Delta }}{}}{\displaystyle \frac{\mathrm{log}p}{p}}+{\displaystyle \underset{\begin{array}{c}p|\mathrm{\Delta }\\ p>\mathrm{log}\mathrm{\Delta }\end{array}}{}}{\displaystyle \frac{\mathrm{log}\mathrm{log}\mathrm{\Delta }}{\mathrm{log}\mathrm{\Delta }}}`$
$`\mathrm{log}\mathrm{log}\mathrm{\Delta }+{\displaystyle \frac{\mathrm{log}\mathrm{log}\mathrm{\Delta }}{\mathrm{log}\mathrm{\Delta }}}{\displaystyle \frac{\mathrm{log}\mathrm{\Delta }}{\mathrm{log}\mathrm{log}\mathrm{\Delta }}}`$
$`\mathrm{log}\mathrm{log}\mathrm{\Delta }+1.`$
Finally, note that $`\mathrm{\Delta }h^{k^2}`$, so that $`\mathrm{log}\mathrm{\Delta }\mathrm{log}h`$. This completes the proof of (2.4).
Now consider $`𝔖()`$. From the definition of $`𝔖()`$, we see that
$$\mathrm{log}𝔖()=\underset{p}{}\left\{\left(\frac{k\nu _p}{p}\right)+O\left(\frac{1}{p^2}\right)\right\}1+\underset{p|\mathrm{\Delta }}{}\frac{1}{p}.$$
The last sum may be bounded in a manner similar to that used for $`\beta ()`$. We have
$`{\displaystyle \underset{p|\mathrm{\Delta }}{}}{\displaystyle \frac{1}{p}}`$ $`{\displaystyle \underset{p\mathrm{log}\mathrm{\Delta }}{}}{\displaystyle \frac{1}{p}}+{\displaystyle \underset{\begin{array}{c}p|\mathrm{\Delta }\\ p>\mathrm{log}\mathrm{\Delta }\end{array}}{}}{\displaystyle \frac{1}{\mathrm{log}\mathrm{\Delta }}}`$
$`\mathrm{log}\mathrm{log}\mathrm{log}\mathrm{\Delta }+{\displaystyle \frac{1}{\mathrm{log}\mathrm{\Delta }}}{\displaystyle \frac{\mathrm{log}\mathrm{\Delta }}{\mathrm{log}\mathrm{log}\mathrm{\Delta }}}`$
$`\mathrm{log}\mathrm{log}\mathrm{log}\mathrm{\Delta }.`$
As noted before, $`\mathrm{log}\mathrm{\Delta }\mathrm{log}h`$. Therefore, there is some constant $`b_k`$ such that $`\mathrm{log}𝔖()b_k\mathrm{log}\mathrm{log}\mathrm{log}h`$, and (2.5) follows. ∎
In our final lemma of this section, we give a computation that was used in (1.10).
###### Lemma 7.
Suppose that $`q`$ is an integer with all of its prime divisors less than $`\sqrt{N}`$. Then there is some absolute constant $`c`$ such that
$$\underset{\begin{array}{c}N<n2N\\ (n,q)=1\end{array}}{}\varpi \varpi (n)=2N\left(\mathrm{log}N+C_0\underset{p|q}{}\frac{\mathrm{log}p}{p}\right)+O(N\mathrm{exp}(c\sqrt{\mathrm{log}N})),$$
where
(2.7)
$$C_0=2\mathrm{log}22\gamma 12\underset{p}{}\frac{\mathrm{log}p}{p(p1}.$$
###### Proof.
We first use the hyperbola method to write
$`{\displaystyle \underset{nx}{}}\varpi \varpi (n)=`$ $`2{\displaystyle \underset{m\sqrt{x}}{}}\varpi (m){\displaystyle \underset{nx/m}{}}\varpi (n)\left({\displaystyle \underset{m\sqrt{x}}{}}\varpi (m)\right)^2`$
$`=`$ $`2x{\displaystyle \underset{p\sqrt{x}}{}}{\displaystyle \frac{\mathrm{log}p}{p}}x+O\left(x\mathrm{exp}(c\sqrt{\mathrm{log}x})\right).`$
Next, we use the classical estimate
$$\underset{px}{}\frac{\mathrm{log}p}{p}=\mathrm{log}x\gamma \underset{p}{}\frac{\mathrm{log}p}{p(p1)}+O(\mathrm{exp}(c\sqrt{\mathrm{log}x})),$$
to get
(2.8)
$$\underset{nx}{}\varpi \varpi (n)=x\mathrm{log}x+C_1x+O(x\mathrm{exp}(c\sqrt{\mathrm{log}x})),$$
where
$$C_1=2\gamma 2\underset{p}{}\frac{\mathrm{log}p}{p(p1)}1.$$
We use (2.8) with $`x=N`$, $`x=2N`$, and take differences to get
(2.9)
$$\underset{N<n2N}{}\varpi \varpi (n)=N\mathrm{log}N+NC_0+O(N\mathrm{exp}(c\sqrt{\mathrm{log}N})).$$
Finally, we note that for a given integer $`q<\sqrt{N}`$,
(2.10) $`{\displaystyle \underset{p|q}{}}{\displaystyle \underset{\begin{array}{c}N<n2N\\ (n,q)=p\end{array}}{}}\varpi \varpi (n)=`$ $`2{\displaystyle \underset{p|q}{}}\mathrm{log}p{\displaystyle \underset{N/p<n2N/p}{}}\varpi (n)`$
$`=`$ $`2N{\displaystyle \underset{p|q}{}}{\displaystyle \frac{\mathrm{log}p}{p}}+O(N\mathrm{exp}(c\sqrt{\mathrm{log}N})).`$
The lemma follows by combining (2.9) and (2.10). ∎
## 3. Proof of Theorem 5
As we noted in the introduction, we take $`\nu _p()`$ to be the number of distinct residue classes mod $`p`$ in $``$. We extend this definition to arbitrary squarefree moduli $`d`$ as follows. Let $`_d`$ be the ring of integers mod $`d`$ and define
(3.1)
$$\mathrm{\Omega }_d()=\{a_d:P(a;)0(modd)\},$$
We define $`\nu _d()`$ to be the cardinality of $`\mathrm{\Omega }_d()`$.
Assume that $`d_1,d_2`$ are squarefree numbers with $`(d_1,d_2)=1`$. The Chinese Remainder Theorem gives an isomorphism
(3.2)
$$\xi :_{d_1}\times _{d_2}_{d_1d_2}.$$
The set $`\mathrm{\Omega }_{d_1d_2}()`$ is the image of $`\mathrm{\Omega }_{d_1}()\times \mathrm{\Omega }_{d_2}()`$ under the isomorphism $`\xi `$, so $`\nu _d()`$ is multiplicative.
Throughout this section, we will take $``$ to be a fixed admissible set, and we will usually write $`\nu _d`$ in place of $`\nu _d()`$.
The left-hand side of (1.25) is
(3.3) $`{\displaystyle \underset{N<n2N}{}}`$ $`\left({\displaystyle \underset{d|P(n;)}{}}\lambda _{d,\mathrm{}_1}\right)\left({\displaystyle \underset{e|P(n;)}{}}\lambda _{e,\mathrm{}_2}\right)`$
$`={\displaystyle \underset{d,e}{}}\lambda _{d,\mathrm{}_1}\lambda _{e,\mathrm{}_2}{\displaystyle \underset{\begin{array}{c}N<n2N\\ [d,e]|P(n;)\end{array}}{}}1`$
$`=N{\displaystyle \underset{d,e}{}}{\displaystyle \frac{\lambda _{d,\mathrm{}_1}\lambda _{e,\mathrm{}_2}}{f([d,e])}}+O\left({\displaystyle \underset{d,e}{}}|\lambda _{d,\mathrm{}_1}\lambda _{e,\mathrm{}_2}r_{[d,e]}|\right)`$
$`=NS_1+O(S_2),`$
say, where
(3.4)
$$f(d)=\frac{d}{\nu _d},$$
and
$$r_d=\underset{\begin{array}{c}N<n2N\\ d|P(n;)\end{array}}{}1\frac{N}{f(d)}.$$
The estimates of $`S_1`$ and $`S_2`$ require the following two lemmas.
###### Lemma 8.
We have
$$\underset{r<R}{}\frac{\mu ^2(r)}{f_1(r)}(\mathrm{log}R/r)^{\mathrm{}}=\frac{\mathrm{}!(\mathrm{log}R)^{k+\mathrm{}}}{𝔖()(k+\mathrm{})!}\left\{1+O(\beta ()𝔖()/\mathrm{log}R)\right\}.$$
###### Proof.
We apply Lemma 5 with
$$\gamma (p)=\nu _p,g(p)=\frac{\nu _p}{p\nu _p}=\frac{1}{f_1(p)}.$$
Now $`\nu _p\mathrm{min}(k,p1),`$ so (2.1) holds with $`A_1=k+1`$. Moreover,
$$\beta ()\underset{wp<z}{}\frac{(\nu _pk)\mathrm{log}p}{p}0$$
and
$$\underset{wp<z}{}\frac{\mathrm{log}p}{p}=\mathrm{log}(z/w)+O(1).$$
Therefore (2.2) holds with $`\kappa =k`$, $`A_2`$ some constant depending only on $`k`$, and
$$L1+\beta ()\beta ().$$
Finally, we note that
$$c_\gamma =\underset{p}{}\left(1\frac{\nu _p}{p}\right)^1\left(1\frac{1}{p}\right)^k=\frac{1}{𝔖()}.$$
###### Lemma 9.
Let $`\lambda _{d,\mathrm{}}`$ be as defined in (1.22). If $`d<R`$ and $`d`$ is squarefree, then
$$|\lambda _{d,\mathrm{}}|(\mathrm{log}R)^{k+\mathrm{}}.$$
###### Proof.
From (1.22), we see that if $`d`$ satisfies the hypotheses of the lemma, then
$`|\lambda _{d,\mathrm{}}|`$ $`={\displaystyle \frac{𝔖()}{\mathrm{}!}}{\displaystyle \frac{f(d)}{f_1(d)}}{\displaystyle \underset{\begin{array}{c}r<R/d\\ (r,d)=1\end{array}}{}}{\displaystyle \frac{\mu ^2(r)}{f_1(r)}}(\mathrm{log}R/rd)^{\mathrm{}}`$
$`={\displaystyle \frac{𝔖()}{\mathrm{}!}}{\displaystyle \underset{t|d}{}}{\displaystyle \frac{1}{f_1(t)}}{\displaystyle \underset{\begin{array}{c}r<R/d\\ (r,d)=1\end{array}}{}}{\displaystyle \frac{\mu ^2(r)}{f_1(r)}}(\mathrm{log}R/rd)^{\mathrm{}}.`$
We move the factor $`1/f_1(t)`$ inside the sum and write $`s=rt`$ to get
$`|\lambda _{d,\mathrm{}}|`$ $`={\displaystyle \frac{𝔖()}{\mathrm{}!}}{\displaystyle \underset{t|d}{}}{\displaystyle \underset{\begin{array}{c}r<R/d\\ (r,d)=1\end{array}}{}}{\displaystyle \frac{\mu ^2(r)}{f_1(rt)}}(\mathrm{log}R/rd)^{\mathrm{}}`$
$`={\displaystyle \frac{𝔖()}{\mathrm{}!}}{\displaystyle \underset{t|d}{}}{\displaystyle \underset{\begin{array}{c}s<Rt/d\\ (s,d)=t\end{array}}{}}{\displaystyle \frac{\mu ^2(s)}{f_1(s)}}(\mathrm{log}Rt/sd)^{\mathrm{}}.`$
For any $`t|d`$, we have $`Rt/d<R`$, so
$$|\lambda _{d,\mathrm{}}|\frac{𝔖()}{\mathrm{}!}(\mathrm{log}R)^{\mathrm{}}\underset{t|d}{}\underset{\begin{array}{c}s<R\\ (s,d)=t\end{array}}{}\frac{\mu ^2(s)}{f_1(s)}.$$
Now for any $`s<R`$, there is a unique $`t|d`$ such that $`(s,d)=t`$. Therefore
$$|\lambda _{d,\mathrm{}}|\frac{𝔖()}{\mathrm{}!}(\mathrm{log}R)^{\mathrm{}}\underset{s<R}{}\frac{\mu ^2(s)}{f_1(s)}.$$
To complete the proof, we observe that
$`{\displaystyle \underset{s<R}{}}{\displaystyle \frac{\mu ^2(s)}{f_1(s)}}`$ $`{\displaystyle \underset{p<R}{}}\left(1+{\displaystyle \frac{1}{f_1(p)}}\right)`$
$`={\displaystyle \underset{p<R}{}}\left(1{\displaystyle \frac{\nu _p}{p}}\right)^1\left(1{\displaystyle \frac{1}{p}}\right)^k{\displaystyle \underset{p<R}{}}\left(1{\displaystyle \frac{1}{p}}\right)^k`$
$`{\displaystyle \frac{(\mathrm{log}R)^k}{𝔖()}}.`$
We now treat $`S_1`$ and $`S_2`$. For $`S_1`$, we begin by writing
$`S_1=`$ $`{\displaystyle \underset{d,e}{}}{\displaystyle \frac{\lambda _{d,\mathrm{}_1}\lambda _{e,\mathrm{}_2}}{f(d)f(e)}}{\displaystyle \underset{\begin{array}{c}r|d\\ r|e\end{array}}{}}f_1(r)`$
$`=`$ $`{\displaystyle \stackrel{\mathrm{}}{}_r}f_1(r)\left({\displaystyle \underset{d}{}}{\displaystyle \frac{\lambda _{dr,\mathrm{}_1}}{f(dr)}}\right)\left({\displaystyle \underset{e}{}}{\displaystyle \frac{\lambda _{er,\mathrm{}_1}}{f(er)}}\right)`$
$`=`$ $`{\displaystyle \stackrel{\mathrm{}}{}_r}{\displaystyle \frac{y_{r,\mathrm{}_1}y_{r,\mathrm{}_2}}{f_1(r)}}`$
$`=`$ $`{\displaystyle \frac{𝔖()^2}{\mathrm{}_1!\mathrm{}_2!}}{\displaystyle \underset{r<R}{}}{\displaystyle \frac{\mu ^2(r)\mathrm{log}^{\mathrm{}_1+\mathrm{}_2}(R/r)}{f_1(r)}}.`$
Lemma 8 now yields the estimate
$$S_1=\left(\genfrac{}{}{0pt}{}{\mathrm{}_1+\mathrm{}_2}{\mathrm{}_1}\right)𝔖()\frac{(\mathrm{log}R)^{k+\mathrm{}_1+\mathrm{}_2}}{(k+\mathrm{}_1+\mathrm{}_2)!}\left\{1+O(\beta ()𝔖()/\mathrm{log}R)\right\}.$$
For $`S_2`$, we first note that
$$|r_d|\nu _dk^{\omega (d)}.$$
We also have the bound for $`\lambda _{d,\mathrm{}}`$ given in Lemma 9. Therefore
$`S_2`$ $`={\displaystyle \underset{d,e<R}{}}|\lambda _{d,\mathrm{}_1}\lambda _{e,\mathrm{}_2}r_{[d,e]}|`$
$`(\mathrm{log}R)^{2k+\mathrm{}_1+\mathrm{}_2}{\displaystyle \stackrel{\mathrm{}}{}_{d,eR}}k^{\omega ([d,e])}`$
$`(\mathrm{log}R)^{4k}{\displaystyle \stackrel{\mathrm{}}{}_{r<R^2}}(3k)^{\omega (r)}.`$
Using Lemma 1, we get
(3.5)
$$S_2R^2(\mathrm{log}R)^{7k}(N/\mathrm{log}N)$$
provided $`R<N^{1/2ϵ}`$.
Theorem 5 follows by combining the above estimates for $`S_1`$ and $`S_2`$.
## 4. Proof of Theorem 6, Part 1
In this section, we consider Theorem 6 under the assumption that $`h_0`$. Our problem is translation invariant in $``$, so we may, without loss of generality, assume that $`h_0=0`$ and $`0`$.
Let $``$ denote the sum on the left-hand side of (1.26). Then
(4.1)
$$=\underset{d,e}{}\lambda _{d,\mathrm{}_1}\lambda _{e,\mathrm{}_2}\underset{\begin{array}{c}N<n2N\\ [d,e]|P(n;)\end{array}}{}\varpi (n)=\underset{d,e}{}\lambda _{d,\mathrm{}_1}\lambda _{e,\mathrm{}_2}\underset{a\mathrm{\Omega }_{[d,e]}()}{}\underset{\begin{array}{c}N<p2N\\ pa(mod[d,e])\end{array}}{}\mathrm{log}p.$$
Now all prime divisors of $`[d,e]`$ are $`<R`$, and $`R<N`$. Therefore, the innermost sum in (4.1) is $`0`$ if $`(a,[d,e])1`$. Accordingly, we need an analogue of $`\mathrm{\Omega }_d()`$ for reduced residue classes. For squarefree $`d`$, we define
(4.2)
$$\mathrm{\Omega }_d^{}()=\{a_d:(a,d)=1\text{ and }P(a;)0(modd)\}.$$
Let $`\nu _d^{}()`$ be the cardinality of $`\mathrm{\Omega }_d^{}()`$. For brevity, we will usually write $`\nu _d^{}`$ in place of $`\nu _d^{}()`$.
When $`d_1,d_2`$ are squarefree and $`(d_1,d_2)=1`$, the set $`\mathrm{\Omega }_{d_1d_2}^{}()`$ is the image of $`\mathrm{\Omega }_{d_1}^{}\times \mathrm{\Omega }_{d_2}^{}`$ under the isomorphism $`\xi `$ of (3.2). Therefore, the function $`\nu ^{}`$ is multiplicative. Moreover, when $`p`$ is prime,
$$\nu _p^{}=\nu _p1,$$
because we are assuming that $`0`$.
In this context, the most natural analogue of $`𝔖()`$ is the product
(4.3)
$$𝔖^{}()=\underset{p}{}\left(1\frac{\nu _p^{}}{p1}\right)\left(1\frac{1}{p}\right)^{k+1}.$$
Note, however that
(4.4) $`𝔖^{}()=`$ $`{\displaystyle \underset{p}{}}\left(1{\displaystyle \frac{\nu _p1}{p1}}\right)\left(1{\displaystyle \frac{1}{p}}\right)^{k+1}`$
$`=`$ $`{\displaystyle \underset{p}{}}\left({\displaystyle \frac{p\nu _p}{p1}}\right)\left({\displaystyle \frac{p1}{p}}\right)\left(1{\displaystyle \frac{1}{p}}\right)^k`$
$`=`$ $`𝔖().`$
Returning to $``$, we write this sum as
(4.5)
$$=\underset{d,e}{}\lambda _{d,\mathrm{}_1}\lambda _{e,\mathrm{}_2}\underset{a\mathrm{\Omega }_{[d,e]}^{}()}{}\underset{\begin{array}{c}N<p2N\\ pa(mod[d,e])\end{array}}{}\mathrm{log}p=NS+O(T),$$
where
(4.6)
$$S=\underset{d,e}{}\frac{\lambda _{d,\mathrm{}_1}\lambda _{e,\mathrm{}_2}\nu _{[d,e]}^{}}{\varphi ([d,e])}$$
and
$$T=\underset{d,e}{}|\lambda _{d,\mathrm{}_1}\lambda _{e,\mathrm{}_2}|\nu _{[d,e]}^{}E^{}(N,[d,e]).$$
By Lemma 9 and Lemma 2,
(4.7)
$$T(\mathrm{log}R)^{2k+\mathrm{}_1+\mathrm{}_2}\stackrel{\mathrm{}}{}_{r<R^2}(3k3)^{\omega (r)}E^{}(N,r)(N/\mathrm{log}N).$$
We now consider the sum $`S`$. We shall define
(4.8)
$$f^{}(r)=\frac{\varphi (r)}{\nu _r^{}}.$$
However, we need to take some care with this definition because there may be terms with $`\nu _r^{}=0`$. However, $`\nu _p^{}=k1`$ for all but finitely many primes $`p`$, so there are at most finitely many primes $`p`$ such that $`\nu _p^{}=0`$. We define
(4.9)
$$A=A()=\underset{\begin{array}{c}p\\ \nu _p^{}()=0\end{array}}{}p,$$
and we use the definition in (4.8) for any $`r`$ with $`(r,A)=1`$. We define $`f_1^{}`$, a function analogous to $`f_1`$, by taking
$$f_1^{}(r)=f^{}\mu (r)$$
for $`r`$ with $`(r,A)=1`$. For future reference, we note that if $`p`$ is a prime and $`pA`$, then
$$f^{}(p)=\frac{p1}{\nu _p1},f_1^{}(p)=\frac{p\nu _p}{\nu _p1}.$$
With this definition of $`f^{}`$, we now have
$$S=\stackrel{}{}_{d,e}\frac{\lambda _{d,\mathrm{}_1}\lambda _{e,\mathrm{}_2}}{f^{}(d)f^{}(e)}\underset{\begin{array}{c}r|d\\ r|e\end{array}}{}f_1^{}(r).$$
Here, and in the sequel, we use $`_{}^{}{}_{}{}^{}`$ to denote that the sum is over values of the indices that are relatively prime to $`A`$. Interchanging the order of summation, we get
(4.10) $`S=`$ $`{\displaystyle \stackrel{}{}_r}f_1^{}(r)\left({\displaystyle \stackrel{}{}_d}{\displaystyle \frac{\lambda _{dr,\mathrm{}_1}}{f^{}(dr)}}\right)\left({\displaystyle \stackrel{}{}_e}{\displaystyle \frac{\lambda _{er,\mathrm{}_2}}{f^{}(er)}}\right)`$
$`=`$ $`{\displaystyle \stackrel{}{}_r}{\displaystyle \frac{y_{r,\mathrm{}_1}^{}y_{r,\mathrm{}_2}^{}}{f_1^{}(r)}},`$
where the quantity $`y_{r,\mathrm{}}^{}`$ is analogous to $`y_{r,\mathrm{}}`$ and is defined as
(4.11)
$$y_{r,\mathrm{}}^{}=\{\begin{array}{cc}\mu (r)f_1^{}(r)\stackrel{}{}_d\frac{\lambda _{dr,\mathrm{}}}{f^{}(dr)}\hfill & \text{ if }(r,A)=1\text{ and }r<R\text{,}\hfill \\ 0\hfill & \text{ otherwise}.\hfill \end{array}$$
Upon using (1.21), the original definition of $`\lambda _{d,\mathrm{}}`$, we see that
$`{\displaystyle \frac{\mu (r)y_{r,\mathrm{}}^{}}{f_1^{}(r)}}`$ $`={\displaystyle \stackrel{}{}_d}{\displaystyle \frac{\lambda _{dr,\mathrm{}}}{f^{}(dr)}}={\displaystyle \stackrel{}{}_d}{\displaystyle \frac{\mu (dr)}{f^{}(dr)}}f(dr){\displaystyle \underset{t}{}}{\displaystyle \frac{y_{rdt,\mathrm{}}}{f_1(rdt)}}`$
$`={\displaystyle \frac{\mu (r)f(r)}{f^{}(r)f_1(r)}}{\displaystyle \stackrel{}{}_{\begin{array}{c}d\\ (d,r)=1\end{array}}}{\displaystyle \frac{\mu (d)f(d)}{f^{}(d)}}{\displaystyle \underset{t}{}}{\displaystyle \frac{y_{rdt,\mathrm{}}}{f_1(dt)}}`$
$`={\displaystyle \frac{\mu (r)f(r)}{f^{}(r)f_1(r)}}{\displaystyle \underset{\begin{array}{c}m\\ (m,r)=1\end{array}}{}}{\displaystyle \frac{y_{rm,\mathrm{}}}{f_1(m)}}{\displaystyle \stackrel{}{}_{d|m}}{\displaystyle \frac{\mu (d)f(d)}{f^{}(d)}}.`$
Note that $`m`$ can be any squarefree integer; we need not have $`(m,A)=1`$. Now
$`{\displaystyle \stackrel{}{}_{d|m}}{\displaystyle \frac{\mu (d)f(d)}{f^{}(d)}}=`$ $`{\displaystyle \underset{p|m,pA}{}}\left(1{\displaystyle \frac{f(p)}{f^{}(p)}}\right)`$
$`=`$ $`{\displaystyle \underset{p|m,pA}{}}\left({\displaystyle \frac{p\nu _p}{\nu _p(p1)}}\right)`$
$`=`$ $`{\displaystyle \underset{p|m}{}}\left({\displaystyle \frac{p\nu _p}{\nu _p(p1)}}\right).`$
We may drop the condition that $`pA`$ in the last line because when $`p|A`$, $`\nu _p=1`$, and $`(p\nu _p)/(\nu _p(p1))=1`$. Therefore
(4.12)
$$\frac{1}{f_1(m)}\stackrel{}{}_{d|m}\frac{\mu (d)f(d)}{f^{}(d)}=\underset{p|m}{}\frac{p\nu _p}{\nu _p(p1)f_1(p)}=\frac{1}{\varphi (m)}.$$
Moreover,
(4.13)
$$\frac{f_1^{}(r)f(r)}{f^{}(r)f_1(r)}=\frac{r}{\varphi (r)}$$
when $`(r,A)=1`$, and so
(4.14)
$$y_{r,\mathrm{}}^{}=\mu ^2(r)\frac{𝔖()}{\mathrm{}!}\frac{r}{\varphi (r)}\underset{\begin{array}{c}m<R/r\\ (m,r)=1\end{array}}{}\frac{\mu ^2(m)}{\varphi (m)}(\mathrm{log}R/rm)^{\mathrm{}}$$
when $`(r,A)=1`$.
For the inner sum, we use Lemma 5 with
$$\gamma (p)=\{\begin{array}{cc}1\hfill & \text{ if }pr\text{,}\hfill \\ 0\hfill & \text{ if }p|r.\hfill \end{array}$$
The hypotheses (2.1) and (2.2) are satisfied with $`\kappa =1`$, some absolute constants $`A_1,A_2`$, and
$$L=\underset{p|r}{}\frac{\mathrm{log}p}{p}+O(1).$$
Let
(4.15)
$$\rho (r)=1+\underset{p|r}{}\frac{\mathrm{log}p}{p},$$
so that $`L\rho (r)`$. With this choice of $`\gamma `$, we have
$$c_\gamma =\underset{p|r}{}\left(1\frac{1}{p}\right)=\frac{\varphi (r)}{r}.$$
We therefore conclude that
(4.16)
$$\underset{\begin{array}{c}m<R/r\\ (m,r)=1\end{array}}{}\frac{\mu ^2(m)}{\varphi (m)}(\mathrm{log}R/rm)^{\mathrm{}}=\frac{\varphi (r)}{r}\frac{(\mathrm{log}R/r)^{\mathrm{}+1}}{\mathrm{}+1}+O\left(\rho (r)(\mathrm{log}2R/r)^{\mathrm{}}\right).$$
We remark parenthetically that Hildebrand gave a more precise formula for this sum in the case $`\mathrm{}=0`$. It is possible to use his result to derive a more accurate version of (4.16), but the above version is sufficient for our purposes.
From (4.16) and (4.14), we deduce that when $`(r,A)=1`$ and $`r<R`$,
(4.17)
$$y_{r,\mathrm{}}^{}=\mu ^2(r)\frac{𝔖()}{(\mathrm{}+1)!}(\mathrm{log}R/r)^{\mathrm{}+1}+O\left(\frac{\mu ^2(r)\rho (r)r}{\varphi (r)}𝔖()(\mathrm{log}2R/r)^{\mathrm{}}\right).$$
We plug this back into our formula for $`S`$ in (4.10) to get
(4.18)
$$S=\stackrel{}{}_{r<R}\frac{y_{r,\mathrm{}_1}^{}y_{r,\mathrm{}_2}^{}}{f_1^{}(r)}=V+O\left(𝔖()^2(\mathrm{log}R)^{\mathrm{}_1+\mathrm{}_2+1}W\right),$$
where
(4.19)
$$V=\frac{𝔖()^2}{(\mathrm{}_1+1)!(\mathrm{}_2+1)!}\stackrel{}{}_{r<R}\frac{\mu ^2(r)}{f_1^{}(r)}(\mathrm{log}R/r)^{\mathrm{}_1+\mathrm{}_2+2}$$
and
(4.20)
$$W=\stackrel{}{}_{r<R}\frac{\mu ^2(r)}{f_1^{}(r)}\frac{\rho (r)r}{\varphi (r)}.$$
We will use Lemma 5 for $`V`$. We will need to estimate a similar sum in Section 6, so it is convenient to have the following lemma that is general enough to cover both situations.
###### Lemma 10.
If $`d`$ is squarefree, $`d<R`$, and $`a`$ is a non-negative integer, then
$`{\displaystyle \stackrel{}{}_{\begin{array}{c}r<R/d\\ (r,d)=1\end{array}}}{\displaystyle \frac{\mu ^2(r)}{f_1^{}(r)}}(\mathrm{log}R/dr)^a=`$ $`{\displaystyle \frac{1}{𝔖()}}{\displaystyle \frac{a!}{(k+a1)!}}(\mathrm{log}R/d)^{k+a1}{\displaystyle \underset{p|d}{}}\left({\displaystyle \frac{p\nu _p}{p1}}\right)`$
$`+O\left((\beta ()+\rho (d))(\mathrm{log}2R/d)^{k+a2}\right).`$
###### Proof.
We apply Lemma 5 with
$$\gamma (p)=\{\begin{array}{cc}\frac{p\nu _p^{}}{p1}\hfill & \text{if }(p,d)=1,\hfill \\ 0\hfill & \text{if }p|d.\hfill \end{array}$$
With this definition for $`\gamma `$, we have
$$g(p)=\frac{\gamma (p)}{p\gamma (p)}=\frac{1}{f_1^{}(p)}$$
when $`(p,Ad)=1`$. Moreover,
$$\nu _p^{}=\nu _p1\mathrm{min}(k1,p2),$$
so (2.1) is true with $`A_1=k`$. For (2.2), we first note that
$`{\displaystyle \underset{wp<z}{}}`$ $`{\displaystyle \frac{\gamma (p)\mathrm{log}p}{p}}={\displaystyle \underset{\begin{array}{c}wp<z\\ (p,d)=1\end{array}}{}}{\displaystyle \frac{(\nu _p1)\mathrm{log}p}{p1}}`$
$`=(k1){\displaystyle \underset{wp<z}{}}{\displaystyle \frac{\mathrm{log}p}{p1}}{\displaystyle \underset{\begin{array}{c}wp<z\\ (p,d)=1\end{array}}{}}{\displaystyle \frac{(k\nu _p)\mathrm{log}p}{p1}}{\displaystyle \underset{\begin{array}{c}wp<z\\ p|d\end{array}}{}}{\displaystyle \frac{(k1)\mathrm{log}p}{p1}}.`$
Now
$`{\displaystyle \underset{wp<z}{}}`$ $`{\displaystyle \frac{\mathrm{log}p}{p1}}=\mathrm{log}(z/w)+O(1),`$
$`{\displaystyle \underset{\begin{array}{c}wp<z\\ (p,d)=1\end{array}}{}}`$ $`{\displaystyle \frac{(k\nu _p)\mathrm{log}p}{p1}}\beta ()+O(1),`$
$`{\displaystyle \underset{\begin{array}{c}wp<z\\ p|d\end{array}}{}}`$ $`{\displaystyle \frac{(k1)\mathrm{log}p}{p1}}(k1)\rho (d),`$
so (2.2) is satisfied with $`\kappa =k1`$, $`A_2`$ some constant depending only on $`k`$, and $`L=\beta ()+(k1)\rho (d)+O(1)\beta ()+\rho (d)`$. Finally, we note that in this situation,
$`c_\gamma `$ $`={\displaystyle \underset{p}{}}\left(1{\displaystyle \frac{\nu _p^{}}{p1}}\right)^1\left(1{\displaystyle \frac{1}{p}}\right)^{k1}{\displaystyle \underset{p|d}{}}\left(1{\displaystyle \frac{\nu _p^{}}{p1}}\right)`$
$`={\displaystyle \frac{1}{𝔖()}}{\displaystyle \underset{p|d}{}}\left({\displaystyle \frac{p\nu _p}{p1}}\right)`$
by (4.4). ∎
From the previous lemma, with $`d=1`$, we see that
(4.21) $`V`$ $`={\displaystyle \frac{𝔖()}{(\mathrm{}_1+1)!(\mathrm{}_2+1)!}}{\displaystyle \frac{(\mathrm{}_1+\mathrm{}_2+2)!}{(k+\mathrm{}_1+\mathrm{}_2+1)!}}(\mathrm{log}R)^{k+\mathrm{}_1+\mathrm{}_2+1}`$
$`+O(\beta ()𝔖()^2(\mathrm{log}R)^{k+\mathrm{}_1+\mathrm{}_2})`$
$`=\left({\displaystyle \genfrac{}{}{0pt}{}{\mathrm{}_1+\mathrm{}_2+2}{\mathrm{}_1+1}}\right)𝔖(){\displaystyle \frac{(\mathrm{log}R)^{k+\mathrm{}_1+\mathrm{}_2+1}}{(k+\mathrm{}_1+\mathrm{}_2+1)!}}+O(\beta ()𝔖()^2(\mathrm{log}R)^{k+\mathrm{}_1+\mathrm{}_2}).`$
The sum $`W`$ may be estimated by relatively trivial means. Now
(4.22) $`W`$ $`={\displaystyle \stackrel{}{}_{r<R}}{\displaystyle \frac{\mu ^2(r)r}{f_1^{}(r)\varphi (r)}}\left(1+{\displaystyle \underset{p|r}{}}{\displaystyle \frac{\mathrm{log}p}{p}}\right)`$
$`={\displaystyle \stackrel{}{}_{r<R}}{\displaystyle \frac{\mu ^2(r)r}{f_1^{}(r)\varphi (r)}}+{\displaystyle \stackrel{}{}_{p<R}}{\displaystyle \frac{\mathrm{log}p}{p}}{\displaystyle \stackrel{}{}_{\begin{array}{c}r<R\\ p|r\end{array}}}{\displaystyle \frac{\mu ^2(r)r}{f_1^{}(r)\varphi (r)}}`$
$`={\displaystyle \stackrel{}{}_{r<R}}{\displaystyle \frac{\mu ^2(r)r}{f_1^{}(r)\varphi (r)}}+{\displaystyle \stackrel{}{}_{p<R}}{\displaystyle \frac{\mathrm{log}p}{f_1^{}(p)\varphi (p)}}{\displaystyle \stackrel{}{}_{\begin{array}{c}r<R/p\\ (r,p)=1\end{array}}}{\displaystyle \frac{\mu ^2(r)r}{f_1^{}(r)\varphi (r)}}`$
$`\left(1+{\displaystyle \stackrel{}{}_{p<R}}{\displaystyle \frac{\mathrm{log}p}{f_1^{}(p)\varphi (p)}}\right)W^{}W^{},`$
where
(4.23)
$$W^{}=\stackrel{}{}_{r<R}\frac{\mu ^2(r)r}{f_1^{}(r)\varphi (r)}=\stackrel{\mathrm{}}{}_{r<R}\frac{\nu _r^{}h(r)}{r},$$
and
(4.24)
$$h(r)=\underset{p|r}{}\frac{p^2}{(p\nu _p)(p1)}.$$
Let $`h_1=h\mu `$, so that
$$h_1(d)=\underset{p|d}{}\frac{p(\nu _p+1)\nu _p}{(p1)(p\nu _p)}.$$
Then
(4.25)
$$W^{}=\stackrel{\mathrm{}}{}_{r<R}\frac{\nu _r^{}}{r}\underset{d|r}{}h_1(d)=\stackrel{\mathrm{}}{}_{d<R}\frac{h_1(d)\nu _d^{}}{d}\stackrel{\mathrm{}}{}_{\begin{array}{c}r<R/d\\ (r,d)=1\end{array}}\frac{\nu _r^{}}{r}\underset{p<R}{}\left(1+\frac{h_1(p)\nu _p^{}}{p}\right)\stackrel{\mathrm{}}{}_{r<R}\frac{\nu _r^{}}{r}.$$
The sum on the right-hand side of (4.25) is $`(\mathrm{log}R)^{k1}`$ by Lemma 1. The product is $`1`$ because
$$\underset{p<R}{}\mathrm{log}\left(1+\frac{h_1(p)\nu _p^{}}{p}\right)\underset{p<R}{}\frac{\nu _p^2}{(p1)(p\nu _p)}1.$$
We conclude that $`W^{}(\mathrm{log}R)^{k1}`$, and so
(4.26)
$$W(\mathrm{log}R)^{k1}.$$
Combining the above with the estimate in (4.21) gives
(4.27)
$$S=\left(\genfrac{}{}{0pt}{}{\mathrm{}_1+\mathrm{}_2+2}{\mathrm{}_1+1}\right)𝔖()\frac{(\mathrm{log}R)^{k+\mathrm{}_1+\mathrm{}_2+1}}{(k+\mathrm{}_1+\mathrm{}_2+1)!}+O\left(\beta ()𝔖()^2(\mathrm{log}R)^{k+\mathrm{}_1+\mathrm{}_2}\right).$$
The first part of Theorem 6 (statement (1.26)) now follows by combining (4.5), (4.7), and (4.27).
## 5. Proof of Theorem 6, Part 2
In this section, we consider Theorem 6 in the case $`h_0`$. As in the previous section, our problem is translation invariant, so we we may assume that $`0`$ and $`^0=\{0\}`$. Consequently, $`P(n;^0)=nP(n;)`$.
Now let $``$ be the left-hand side of (1.27). If $`n`$ is a prime with $`N<n2N`$, then $`1`$ is the only divisor of $`n`$ less than $`N`$. When $`d<R<N`$, we have $`d|P(n;)`$ if and only if $`d|P(n;^0)`$. Consequently,
(5.1)
$$=\underset{d,e}{}\lambda _{d,\mathrm{}_1}\lambda _{e,\mathrm{}_2}\underset{\begin{array}{c}N<n2N\\ [d,e]|P(n;^0)\end{array}}{}\varpi (n)=\underset{d,e}{}\lambda _{d,\mathrm{}_1}\lambda _{e,\mathrm{}_2}\underset{a\mathrm{\Omega }_{[d,e]}^{}(^0)}{}\underset{\begin{array}{c}N<p2N\\ pa(mod[d,e])\end{array}}{}\mathrm{log}p.$$
In parallel to the argument in (4.5) through (4.7), we find that
$$=NS+T,$$
where
(5.2)
$$S=\underset{d,e}{}\frac{\lambda _{d,\mathrm{}_1}\lambda _{e,\mathrm{}_2}\nu _{[d,e]}^{}(^0)}{\varphi ([d,e])}$$
and
$$T=\underset{d,e}{}|\lambda _{d,\mathrm{}_1}\lambda _{e,\mathrm{}_2}|\nu _{[d,e]}^{}(^0)E^{}(N,[d,e])N/\mathrm{log}N.$$
Therefore
(5.3)
$$=NS+O(N/\mathrm{log}N).$$
The rest of this section is devoted to evaluating the sum $`S`$.
For brevity, we write $`\nu _r^{}`$ for $`\nu _r^{}(^0)`$. Let
$$A_0=A(^0)=\underset{\begin{array}{c}p\\ \nu _p^{}=0\end{array}}{}p.$$
For squarefree $`r`$ with $`(r,A_0)=1`$, we define
(5.4)
$$f^{}(r)=\frac{\varphi (r)}{\nu _r^{}}=\underset{p|r}{}\left(\frac{p1}{\nu _p^{}}\right),$$
and
(5.5)
$$f_1^{}(r)=f^{}\mu (r)=\underset{p|r}{}\left(\frac{p1\nu _p^{}}{\nu _p^{}}\right)$$
Note that
$$\nu _p^{}=\{\begin{array}{cc}\nu _p\hfill & \text{if }0\mathrm{\Omega }_p()\hfill \\ \nu _p1\hfill & \text{if }0\mathrm{\Omega }_p().\hfill \end{array}$$
We are assuming that $`0`$, so there are only finitely many primes $`p`$ with $`0\mathrm{\Omega }_p()`$. Let
(5.6)
$$B_0=B_0()=\underset{\begin{array}{c}p\\ \nu _p^{}=\nu _p1\end{array}}{}p=\underset{\begin{array}{c}p\\ 0\mathrm{\Omega }_p()\end{array}}{}p.$$
In fact, $`0\mathrm{\Omega }_p()`$ if and only if $`p`$ divides $`h`$ for some $`h`$. Therefore $`B_0`$ is the squarefree kernel of the product of all elements of $``$.
For future reference, we note that when $`(r,A_0)=1`$,
$$f^{}(r)=\underset{\begin{array}{c}p|r\\ pB_0\end{array}}{}\left(\frac{p1}{\nu _p}\right)\underset{\begin{array}{c}p|r\\ p|B_0\end{array}}{}\left(\frac{p1}{\nu _p1}\right)$$
and
$$f_1^{}(r)=\underset{\begin{array}{c}p|r\\ pB_0\end{array}}{}\left(\frac{p1\nu _p}{\nu _p}\right)\underset{\begin{array}{c}p|r\\ p|B_0\end{array}}{}\left(\frac{p\nu _p}{\nu _p1}\right).$$
With the above definitions of $`f^{}`$ and $`f_1^{}`$, we may write
$`S=`$ $`{\displaystyle \stackrel{}{}_{d,e}}{\displaystyle \frac{\lambda _{d,\mathrm{}_1}\lambda _{e,\mathrm{}_2}}{f^{}([d,e])}}`$
$`=`$ $`{\displaystyle \stackrel{}{}_{d,e}}{\displaystyle \frac{\lambda _{d,\mathrm{}_1}\lambda _{e,\mathrm{}_2}}{f^{}(d)f^{}(e)}}{\displaystyle \underset{\begin{array}{c}r|d\\ r|e\end{array}}{}}f_1^{}(r)`$
$`=`$ $`{\displaystyle \stackrel{}{}_r}f_1^{}(r)\left({\displaystyle \stackrel{}{}_d}{\displaystyle \frac{\lambda _{dr,\mathrm{}_1}}{f^{}(dr)}}\right)\left({\displaystyle \stackrel{}{}_e}{\displaystyle \frac{\lambda _{er,\mathrm{}_2}}{f^{}(er)}}\right),`$
where $`_{}^{}{}_{}{}^{}`$ denotes that the sum is over values of the indices that are relatively prime to $`A_0`$. We get
(5.7)
$$S=\stackrel{}{}_r\frac{y_{r,\mathrm{}_1}^{}y_{r,\mathrm{}_2}^{}}{f_1^{}(r)},$$
where we define
(5.8)
$$y_{r,\mathrm{}}^{}=\{\begin{array}{cc}\mu (r)f_1^{}(r)\stackrel{}{}_d\frac{\lambda _{dr,\mathrm{}}}{f^{}(dr)}\hfill & \text{ if }(r,A_0)=1\text{ and }r<R\text{,}\hfill \\ 0\hfill & \text{ otherwise.}\hfill \end{array}$$
Upon using (1.21), our original definition of $`\lambda _{d,\mathrm{}}`$, we see that
$`{\displaystyle \frac{\mu (r)y_{r,\mathrm{}}^{}}{f_1^{}(r)}}`$ $`={\displaystyle \stackrel{}{}_d}{\displaystyle \frac{\lambda _{dr,\mathrm{}}}{f^{}(dr)}}={\displaystyle \stackrel{}{}_d}{\displaystyle \frac{\mu (dr)}{f^{}(dr)}}f(dr){\displaystyle \underset{t}{}}{\displaystyle \frac{y_{rdt,\mathrm{}}}{f_1(rdt)}}`$
$`={\displaystyle \frac{\mu (r)f(r)}{f^{}(r)f_1(r)}}{\displaystyle \stackrel{}{}_{\begin{array}{c}d\\ (d,r)=1\end{array}}}{\displaystyle \frac{\mu (d)f(d)}{f^{}(d)}}{\displaystyle \underset{t}{}}{\displaystyle \frac{y_{rdt,\mathrm{}}}{f_1(dt)}}`$
$`={\displaystyle \frac{\mu (r)f(r)}{f^{}(r)f_1(r)}}{\displaystyle \underset{\begin{array}{c}m\\ (m,r)=1\end{array}}{}}{\displaystyle \frac{y_{rm,\mathrm{}}}{f_1(m)}}{\displaystyle \stackrel{}{}_{d|m}}{\displaystyle \frac{\mu (d)f(d)}{f^{}(d)}}.`$
Now
$$\stackrel{}{}_{d|m}\frac{\mu (d)f(d)}{f^{}(d)}=\underset{\begin{array}{c}p|m\\ pA_0\end{array}}{}\left(1\frac{f(p)}{f^{}(p)}\right)=\underset{p|m}{}\left(1\frac{p\nu _p^{}}{(p1)\nu _p}\right).$$
The condition $`pA_0`$ can be dropped because $`\nu _p^{}=0`$ when $`p|A_0`$. Therefore
$`{\displaystyle \stackrel{}{}_{d|m}}{\displaystyle \frac{\mu (d)f(d)}{f^{}(d)}}=`$ $`{\displaystyle \underset{\begin{array}{c}p|m\\ pB_0\end{array}}{}}\left(1{\displaystyle \frac{p}{p1}}\right){\displaystyle \underset{\begin{array}{c}p|m\\ p|B_0\end{array}}{}}\left(1{\displaystyle \frac{p(\nu _p1)}{\nu _p(p1)}}\right)`$
$`=`$ $`{\displaystyle \frac{\mu (m)}{\varphi (m)}}f_2(m),`$
where $`f_2`$ is the multiplicative function defined by
(5.9)
$$f_2(p)=\{\begin{array}{cc}1\hfill & \text{ if }pB_0,\hfill \\ f_1(p)\hfill & \text{ if }p|B_0.\hfill \end{array}$$
In other words,
$$f_2(m)=\mu ((m,B_0))f_1((m,B_0)).$$
Therefore
(5.10) $`y_{r,\mathrm{}}^{}=`$ $`\mu ^2(r){\displaystyle \frac{f_1^{}(r)f(r)}{f^{}(r)f_1(r)}}{\displaystyle \underset{\begin{array}{c}m<R/r\\ (m,r)=1\end{array}}{}}{\displaystyle \frac{y_{rm,\mathrm{}}}{f_1(m)}}{\displaystyle \frac{\mu (m)}{\varphi (m)}}f_2(m)`$
$`=`$ $`\mu ^2(r){\displaystyle \frac{𝔖()}{\mathrm{}!}}{\displaystyle \frac{f_1^{}(r)f(r)}{f^{}(r)f_1(r)}}{\displaystyle \underset{\begin{array}{c}m<R/r\\ (m,r)=1\end{array}}{}}{\displaystyle \frac{\mu (m)f_2(m)(\mathrm{log}R/rm)^{\mathrm{}}}{f_1(m)\varphi (m)}}.`$
The sum
$$\underset{\begin{array}{c}m=1\\ (m,r)=1\end{array}}{\overset{\mathrm{}}{}}\frac{\mu (m)f_2(m)}{f_1(m)\varphi (m)}$$
converges, and so one would expect that
$$y_{r,\mathrm{}}^{}\mu ^2(r)\frac{𝔖()}{\mathrm{}!}(\mathrm{log}R/r)^{\mathrm{}}\frac{f_1^{}(r)f(r)}{f^{}(r)f_1(r)}\underset{\begin{array}{c}m=1\\ (m,r)=1\end{array}}{\overset{\mathrm{}}{}}\frac{\mu (m)f_2(m)}{f_1(m)\varphi (m)}$$
when $`r<R`$ and $`(r,A_0)=1`$. From Lemma 11 below, we would then obtain
$$y_{r,\mathrm{}}^{}\mu ^2(r)\frac{𝔖(^0)}{\mathrm{}!}(\mathrm{log}R/r)^{\mathrm{}},$$
and we will ultimately prove this. This asymptotic relation should be compared to (1.20) and (4.17).
###### Lemma 11.
If $`r`$ is squarefree and $`(r,A_0)=1`$, then
$$\frac{f_1^{}(r)f(r)}{f^{}(r)f_1(r)}\underset{\begin{array}{c}m=1\\ (m,r)=1\end{array}}{\overset{\mathrm{}}{}}\frac{\mu (m)f_2(m)}{f_1(m)\varphi (m)}=\frac{𝔖(^0)}{𝔖()}.$$
###### Proof.
For $`r`$ satisfying our hypotheses, it is convenient to define
(5.11)
$$F(r)=\frac{f_1^{}(r)f(r)}{f^{}(r)f_1(r)}\text{ and }G(r)=\underset{\begin{array}{c}m=1\\ (m,r)=1\end{array}}{\overset{\mathrm{}}{}}\frac{\mu (m)f_2(m)}{f_1(m)\varphi (m)},$$
so that the left-hand side of the proposed result is $`F(r)G(r)`$. We begin by noting that
$$F(r)=\underset{p|r}{}F(p)=\underset{p|r}{}\frac{p(p1\nu _p^{})}{(p1)(p\nu _p)}.$$
Moreover,
$`G(r)=`$ $`{\displaystyle \underset{pr}{}}\left(1{\displaystyle \frac{f_2(p)}{\varphi (p)f_1(p)}}\right)`$
$`=`$ $`{\displaystyle \underset{\begin{array}{c}pB_0\\ pr\end{array}}{}}{\displaystyle \frac{p(p1\nu _p)}{(p1)(p\nu _p)}}{\displaystyle \underset{\begin{array}{c}p|B_0\\ pr\end{array}}{}}{\displaystyle \frac{p}{p1}}`$
$`=`$ $`{\displaystyle \underset{pr}{}}{\displaystyle \frac{p(p1\nu _p^{})}{(p1)(p\nu _p)}}={\displaystyle \underset{pr}{}}F(p).`$
In the last line, we used the fact that $`\nu _p^{}=\nu _p`$ if $`pB_0`$ and $`\nu _p^{}=\nu _p1`$ if $`pB_0`$. Combining the last two results yields
(5.12)
$$F(r)G(r)=\underset{p}{}\frac{p(p1\nu _p^{})}{(p1)(p\nu _p)}=\underset{p}{}F(p).$$
On the other hand, if we replace $``$ by $`^0`$ and $`k`$ by $`k+1`$ in (4.4), then we obtain
$$𝔖(^0)=𝔖^{}(^0)=\underset{p}{}\left(1\frac{\nu _p^{}}{p1}\right)\left(1\frac{1}{p}\right)^k.$$
We combine this with the definition of $`𝔖()`$ given in (1.1) to get
(5.13)
$$\frac{𝔖(^0)}{𝔖()}=\underset{p}{}\frac{p(p1\nu _p^{})}{(p1)(p\nu _p)}=\underset{p}{}F(p).$$
The lemma follows by comparing this with (5.12). ∎
###### Lemma 12.
Suppose $`\mathrm{}1`$. If $`r<R`$ and $`(r,A_0)=1`$, then
(5.14)
$$y_{r,\mathrm{}}^{}=\mu ^2(r)\frac{𝔖(^0)}{\mathrm{}!}(\mathrm{log}R/r)^{\mathrm{}}+O\left(\mu ^2(r)\beta (^0)𝔖(^0)(\mathrm{log}2R/r)^\mathrm{}1\right).$$
###### Proof.
From the definition of $`y_{r,\mathrm{}}^{}`$ in (5.8), the lemma is trivial if $`r`$ is not squarefree. For the remainder of the proof, we assume that $`r`$ is squarefree, $`(r,A_0)=1`$, and $`r<R`$.
We start from the expression for $`y_{r,\mathrm{}}^{}`$ given in (5.10). For a given $`m`$ in the inner sum, write $`m=\delta n`$, where $`\delta |B_0`$ and $`(n,B_0)=1`$. Then $`f_2(m)=\mu (\delta )f_1(\delta )`$ and
$$\frac{\mu (m)f_2(m)}{f_1(m)\varphi (m)}=\frac{\mu ^2(\delta )\mu (n)}{\varphi (\delta )\varphi (n)f_1(n)}.$$
Therefore (5.10) may be transformed into
$$y_{r,\mathrm{}}^{}=\frac{𝔖()F(r)}{\mathrm{}!}\underset{\begin{array}{c}\delta |B_0\\ (\delta ,r)=1\end{array}}{}\frac{\mu ^2(\delta )}{\varphi (\delta )}\underset{\begin{array}{c}n<R/r\delta \\ (n,rB_0)=1\end{array}}{}\frac{\mu (n)}{f_1(n)\varphi (n)}(\mathrm{log}R/r\delta n)^{\mathrm{}}.$$
If we set
(5.15)
$$B_1=\underset{\begin{array}{c}p|B_0\\ pr\end{array}}{}p=\frac{B_0}{(B_0,r)},$$
then the above equation for $`y_{r,\mathrm{}}^{}`$ may be written as
(5.16)
$$y_{r,\mathrm{}}^{}=\frac{𝔖()F(r)}{\mathrm{}!}\underset{\delta |B_1}{}\frac{\mu ^2(\delta )}{\varphi (\delta )}\underset{\begin{array}{c}n<R/r\delta \\ (n,rB_1)=1\end{array}}{}\frac{\mu (n)}{f_1(n)\varphi (n)}(\mathrm{log}R/r\delta n)^{\mathrm{}}.$$
For future reference, note that $`B_0|rB_1`$.
Now let
(5.17)
$$Y(x;d,\mathrm{})=\underset{\begin{array}{c}n<x\\ (n,d)=1\end{array}}{}\frac{\mu (n)}{f_1(n)\varphi (n)}(\mathrm{log}x/n)^{\mathrm{}},$$
so that the innermost sum in (5.16) is $`Y(R/r\delta ;rB_1,\mathrm{})`$.
Now assume that $`\mathrm{}1`$. We begin our analysis of $`Y`$ by writing
(5.18) $`Y(x;d,\mathrm{})=`$ $`{\displaystyle \underset{\begin{array}{c}n<x\\ (n,d)=1\end{array}}{}}{\displaystyle \frac{\mu (n)}{f_1(n)\varphi (n)}}{\displaystyle _n^x}\mathrm{}(\mathrm{log}x/u)^\mathrm{}1{\displaystyle \frac{du}{u}}`$
$`=`$ $`{\displaystyle _1^x}{\displaystyle \frac{\mathrm{}(\mathrm{log}x/u)^\mathrm{}1}{u}}{\displaystyle \underset{\begin{array}{c}n<u\\ (n,d)=1\end{array}}{}}{\displaystyle \frac{\mu (n)}{f_1(n)\varphi (n)}}du`$
$`=`$ $`Y_1(x;d,\mathrm{})Y_2(x;d,\mathrm{}),`$
where
(5.19)
$$Y_1(x;d,\mathrm{})=_1^x\frac{\mathrm{}(\mathrm{log}x/u)^\mathrm{}1}{u}\underset{\begin{array}{c}n=1\\ (n,d)=1\end{array}}{\overset{\mathrm{}}{}}\frac{\mu (n)}{f_1(n)\varphi (n)}du,$$
and
(5.20)
$$Y_2(x;d,\mathrm{})=_1^x\frac{\mathrm{}(\mathrm{log}x/u)^\mathrm{}1}{u}\underset{\begin{array}{c}nu\\ (n,d)=1\end{array}}{}\frac{\mu (n)}{f_1(n)\varphi (n)}du.$$
We see immediately that
$$Y_1(x;d,\mathrm{})=(\mathrm{log}x)^{\mathrm{}}\underset{pd}{}\left(1\frac{1}{f_1(p)\varphi (p)}\right).$$
If we assume that $`B_0|d`$, then we may write
(5.21)
$$Y_1(x;d,\mathrm{})=(\mathrm{log}x)^{\mathrm{}}\underset{pd}{}F(p).$$
For $`Y_2(x;d,\mathrm{})`$ we bound the sum inside the integrand as
(5.22)
$$\left|\underset{\begin{array}{c}nu\\ (n,d)=1\end{array}}{}\frac{\mu (n)}{f_1(n)\varphi (n)}\right|\underset{nu}{}\frac{\mu ^2(n)}{f_1(n)\varphi (n)}=_u^{\mathrm{}}\left(\underset{un<v}{}\frac{\mu ^2(n)n}{f_1(n)\varphi (n)}\right)\frac{dv}{v^2}.$$
Now let
(5.23)
$$W^{}(v)=\underset{n<v}{}\frac{\mu ^2(n)n}{f_1(n)\varphi (n)}.$$
This sum is very similar to the sum $`W^{}`$ defined in (4.23); in fact,
$$W^{}(v)=\stackrel{\mathrm{}}{}_{n<v}\frac{\nu _nh(n)}{n},$$
where $`h`$ was defined in (4.24). We have, similarly to (4.25),
$$W^{}(v)=\stackrel{\mathrm{}}{}_{n<v}\frac{\nu _n}{n}\underset{d|n}{}h_1(d)=\stackrel{\mathrm{}}{}_{d<v}\frac{h_1(d)\nu _d}{d}\stackrel{\mathrm{}}{}_{\begin{array}{c}n<v/d\\ (n,d)=1\end{array}}\frac{\nu _n}{n}\underset{p<n}{}\left(1+\frac{h_1(p)\nu _p}{p}\right)\stackrel{\mathrm{}}{}_{n<v}\frac{\nu _n}{n}.$$
The sum on the right-hand side is $`(\mathrm{log}2v)^k`$ by Lemma 1. The product on the right hand side is $`1`$ because
$$\underset{p<v}{}\mathrm{log}\left(1+\frac{h_1(p)\nu _p}{p}\right)\underset{p<v}{}\frac{\nu _p^2}{(p1)(p\nu _p)}1.$$
Therefore
(5.24)
$$W^{}(v)(\mathrm{log}2v)^k.$$
Now we use (5.24) in (5.22) to get
$$\left|\underset{\begin{array}{c}nu\\ (n,d)=1\end{array}}{}\frac{\mu (n)}{f_1(n)\varphi (n)}\right|_u^{\mathrm{}}\frac{(\mathrm{log}2v)^k}{v^2}𝑑v\frac{(\mathrm{log}2u)^k}{u}.$$
We use this in (5.20) to get
(5.25)
$$Y_2(x;d,\mathrm{})(\mathrm{log}2x)^\mathrm{}1_1^x(\mathrm{log}2v)^k\frac{dv}{v^2}(\mathrm{log}2x)^\mathrm{}1.$$
Combining this with (5.21) gives
(5.26)
$$Y(x;d,\mathrm{})=(\mathrm{log}x)^{\mathrm{}}\underset{pd}{}F(p)+O((\mathrm{log}2x)^\mathrm{}1)$$
when $`B_0|d`$.
Now we use (5.26) with $`d=rB_1`$ in (5.16) to obtain
(5.27) $`y_{r,\mathrm{}}^{}={\displaystyle \frac{𝔖()}{\mathrm{}!}}`$ $`\left({\displaystyle \underset{pB_1}{}}F(p)\right){\displaystyle \underset{\begin{array}{c}\delta |B_1\\ \delta <R/r\end{array}}{}}{\displaystyle \frac{\mu ^2(\delta )}{\varphi (\delta )}}(\mathrm{log}R/r\delta )^{\mathrm{}}`$
$`+O\left(𝔖()F(r){\displaystyle \underset{\begin{array}{c}\delta |B_1\\ \delta <R/r\end{array}}{}}{\displaystyle \frac{\mu ^2(\delta )}{\varphi (\delta )}}(\mathrm{log}2R/r\delta )^\mathrm{}1\right).`$
The error term in (5.27) is
$``$ $`𝔖()F(r)(\mathrm{log}2R/r)^\mathrm{}1{\displaystyle \underset{\delta |B_1}{}}{\displaystyle \frac{\mu ^2(\delta )}{\varphi (\delta )}}`$
$``$ $`𝔖()\left({\displaystyle \underset{p|rB_1}{}}F(p)\right)(\mathrm{log}2R/r)^\mathrm{}1`$
$``$ $`𝔖(^0)(\mathrm{log}2R/r)^\mathrm{}1\left({\displaystyle \underset{prB_1}{}}F(p)\right)^1.`$
We have used (5.13) in the last line. Now when $`pB_0`$,
$$F(p)^1=\left(1\frac{\nu _p}{(p1)(p\nu _p)}\right)^1=1+O(1/p^2),$$
so
$$\left(\underset{prB_1}{}F(p)\right)^11.$$
Therefore the error term in (5.27) is
(5.28)
$$𝔖(^0)(\mathrm{log}2R/r)^\mathrm{}1.$$
Now we consider the main term in (5.27), which we write as
(5.29)
$$\frac{𝔖()}{\mathrm{}!}\left(\underset{pB_1}{}F(p)\right)\left\{M_1M_2M_3\right\},$$
where
$`M_1=`$ $`(\mathrm{log}R/r)^{\mathrm{}}{\displaystyle \underset{\delta |B_1}{}}{\displaystyle \frac{\mu ^2(\delta )}{\varphi (\delta )}},`$
$`M_2=`$ $`(\mathrm{log}R/r)^{\mathrm{}}{\displaystyle \underset{\begin{array}{c}\delta |B_1\\ \delta R/r\end{array}}{}}{\displaystyle \frac{\mu ^2(\delta )}{\varphi (\delta )}},`$
$`M_3=`$ $`{\displaystyle \underset{\begin{array}{c}\delta |B_1\\ \delta <R/r\end{array}}{}}{\displaystyle \frac{\mu ^2(\delta )}{\varphi (\delta )}}\left\{(\mathrm{log}R/r)^{\mathrm{}}(\mathrm{log}R/r\delta )^{\mathrm{}}\right\}.`$
For $`M_1`$, we note that
$$\underset{\delta |B_1}{}\frac{\mu ^2(\delta )}{\varphi (\delta )}=\underset{p|B_1}{}\frac{p}{p1}=\underset{p|B_1}{}F(p).$$
Therefore
(5.30)
$$\frac{𝔖()}{\mathrm{}!}M_1\underset{pB_1}{}F(p)=\frac{𝔖()}{\mathrm{}!}(\mathrm{log}R/r)^{\mathrm{}}\underset{p}{}F(p)=\frac{𝔖(^0)}{\mathrm{}!}(\mathrm{log}R/r)^{\mathrm{}}.$$
by (5.13).
For $`M_2`$, we note that
$$\underset{\begin{array}{c}\delta |B_1\\ \delta R/r\end{array}}{}\frac{\mu ^2(\delta )}{\varphi (\delta )}\underset{\delta |B_1}{}\frac{\mu ^2(\delta )}{\varphi (\delta )}\frac{\mathrm{log}\delta }{(\mathrm{log}2R/r)},$$
and
$`{\displaystyle \underset{\delta |B_1}{}}{\displaystyle \frac{\mu ^2(\delta )}{\varphi (\delta )}}\mathrm{log}\delta =`$ $`{\displaystyle \underset{\delta |B_1}{}}{\displaystyle \frac{\mu ^2(\delta )}{\varphi (\delta )}}{\displaystyle \underset{p|\delta }{}}\mathrm{log}p={\displaystyle \underset{p|B_1}{}}{\displaystyle \frac{\mathrm{log}p}{p1}}{\displaystyle \underset{\delta |B_1/p}{}}{\displaystyle \frac{\mu ^2(\delta )}{\varphi (\delta )}}`$
$`=`$ $`{\displaystyle \underset{p|B_1}{}}{\displaystyle \frac{\mathrm{log}p}{p1}}{\displaystyle \frac{B_1/p}{\varphi (B_1/p)}}={\displaystyle \frac{B_1}{\varphi (B_1)}}{\displaystyle \underset{p|B_1}{}}{\displaystyle \frac{\mathrm{log}p}{p}}`$
$`=`$ $`F(B_1){\displaystyle \underset{p|B_1}{}}{\displaystyle \frac{\mathrm{log}p}{p}}.`$
Now if $`p|B_1`$, then $`p|B_0`$ and $`\nu _p(^0)k.`$ Therefore
$$\underset{p|B_1}{}\frac{\mathrm{log}p}{p}\underset{p}{}\frac{(k+1\nu _p(^0))\mathrm{log}p}{p}=\beta (^0).$$
Consequently,
(5.31)
$$\underset{\delta |B_1}{}\frac{\mu ^2(\delta )}{\varphi (\delta )}\mathrm{log}\delta F(B_1)\beta (^0),$$
and so
(5.32) $`{\displaystyle \frac{𝔖()}{\mathrm{}!}}M_2{\displaystyle \underset{pB_1}{}}F(p)`$ $`(\mathrm{log}2R/r)^\mathrm{}1𝔖()\beta (^0)F(B_1){\displaystyle \underset{pB_1}{}}F(p)`$
$``$ $`𝔖(^0)\beta (^0)(\mathrm{log}2R/r)^\mathrm{}1.`$
For $`M_3`$, we note that when $`\delta R/r`$,
$`(\mathrm{log}R/r)^{\mathrm{}}`$ $`(\mathrm{log}R/r\delta )^{\mathrm{}}`$
$`=`$ $`(\mathrm{log}\delta )\left\{(\mathrm{log}R/r)^\mathrm{}1+(\mathrm{log}R/r\delta )(\mathrm{log}R/r)^\mathrm{}2+\mathrm{}+(\mathrm{log}R/r\delta )^\mathrm{}1\right\}`$
$``$ $`(\mathrm{log}\delta )(\mathrm{log}R/r)^\mathrm{}1.`$
Thus
$$M_3(\mathrm{log}2R/r)^\mathrm{}1\underset{\delta |B_1}{}\frac{\mu ^2(\delta )}{\varphi (\delta )}\mathrm{log}\delta (\mathrm{log}2R/r)^\mathrm{}1F(B_1)\beta (^0)$$
by (5.31), and so
(5.33)
$$\frac{𝔖()}{\mathrm{}!}M_3\underset{pB_1}{}F(p)(\mathrm{log}2R/r)^\mathrm{}1𝔖(^0)\beta (^0).$$
Combining the estimates (5.28),(5.30), (5.32), and (5.33) gives the proof of Lemma 12. ∎
In reference to the above lemma, we remark that with a bit more work we could give an estimate valid for $`y_{r,0}`$ with a somewhat weaker error term. However, we omit this because it is not necessary for the proof of Theorem 1.
We can now complete the estimate of $`S`$. From (5.2) and Lemma 12, we see that
(5.34)
$$S=V^{}+O\left(𝔖(^0)^2\beta (^0)(\mathrm{log}R)^{\mathrm{}_1+\mathrm{}_21}W^{}\right),$$
where
$$V^{}=\frac{𝔖(^0)^2}{\mathrm{}_1!\mathrm{}_2!}\stackrel{}{}_{r<R}\frac{\mu ^2(r)}{f_1^{}(r)}(\mathrm{log}R/r)^{\mathrm{}_1+\mathrm{}_2},$$
and
$$W^{}=\stackrel{}{}_{r<R}\frac{\mu ^2(r)}{f_1^{}(r)}.$$
Now $`V^{}`$ is the same as the sum $`V`$ in (4.19) except that $``$ has been replaced by $`^0`$, $`k`$ has been replaced by $`k+1`$, and $`\mathrm{}_1,\mathrm{}_2`$ have been replaced by $`\mathrm{}_11,\mathrm{}_21`$ respectively. From (4.21), we see that
(5.35)
$$V^{}=\left(\genfrac{}{}{0pt}{}{\mathrm{}_1+\mathrm{}_2}{\mathrm{}_1}\right)𝔖(^0)\frac{(\mathrm{log}R)^{k+\mathrm{}_1+\mathrm{}_2}}{(k+\mathrm{}_1+\mathrm{}_2)!}+O(\beta (^0)𝔖(^0)^2(\mathrm{log}R)^{k+\mathrm{}_1+\mathrm{}_21}).$$
For $`W^{}`$, we use Lemma 10 with $`a=0`$, $`d=1`$, $`f^{}`$ replaced by $`f^{}`$, $`k`$ replaced by $`k+1`$ to get
(5.36)
$$W^{}(\mathrm{log}R)^k.$$
Now we combine (5.34), (5.35), and (5.36) to get
(5.37)
$$S=\left(\genfrac{}{}{0pt}{}{\mathrm{}_1+\mathrm{}_2}{\mathrm{}_1}\right)𝔖(^0)\frac{(\mathrm{log}R)^{k+\mathrm{}_1+\mathrm{}_2}}{(k+\mathrm{}_1+\mathrm{}_2)!}+O(\beta (^0)𝔖(^0)^2(\mathrm{log}R)^{k+\mathrm{}_1+\mathrm{}_21}).$$
Equation (1.27) now follows by combining this with (5.3).
## 6. Proof of Theorem 7
We may again assume, without loss of generality, that $`h_0=0`$. Accordingly, we assume throughout this section that $`0`$.
Let $``$ denote the sum on the left-hand side in the statement of Theorem 7. Then
(6.1)
$$=\underset{d,e}{}\lambda _{d,\mathrm{}_1}\lambda _{e,\mathrm{}_2}\underset{a\mathrm{\Omega }_{[d,e]}()}{}\underset{\begin{array}{c}N<n2N\\ na(mod[d,e])\end{array}}{}\varpi \varpi (n).$$
In this sum, we have $`d,e<R<\sqrt{N}`$, so $`[d,e]`$ has no prime divisors exceeding $`\sqrt{N}`$. On the other hand, if $`N<n2N`$ and $`\varpi \varpi (n)>0`$, then $`n`$ is a product of two primes, at least one of which must exceed $`\sqrt{N}`$. Therefore, the inner sum in (6.1) will be 0 unless $`(a,[d,e])=1`$ or $`(a,[d,e])=p`$ for some prime $`p<R`$.
We write
$$=_1+_2,$$
where $`_1`$ is the sum in (6.1) with the extra condition that $`(a,[d,e])=1`$, and $`_2`$ is the sum in (6.1) with the extra condition that $`(a,[d,e])=p`$ for some prime $`p`$.
Before analyzing $`_2`$, it is useful to note that when $`r`$ is squarefree and $`(a,r)=p`$,
$`{\displaystyle \underset{\begin{array}{c}N<n2N\\ na(modr)\end{array}}{}}\varpi \varpi (n)=`$ $`2\mathrm{log}p{\displaystyle \underset{\begin{array}{c}\frac{N}{p}<m\frac{2N}{p}\\ m\frac{a}{p}(mod\frac{r}{p})\end{array}}{}}\varpi (m)`$
$`=`$ $`{\displaystyle \frac{2N}{\varphi (r)}}{\displaystyle \frac{(\mathrm{log}p)\varphi (p)}{p}}+O(E^{}(N/p,r/p)).`$
When $`r`$ is squarefree and $`p`$ is a prime dividing $`r`$, we define
(6.2)
$$\mathrm{\Omega }_{r,p}^{}()=\{a_r:(a,r)=p\text{ and }P(a;)0(modr)\}.$$
Let $`\nu _{r,p}^{}=\nu _{r,p}^{}()`$ be the cardinality of $`\mathrm{\Omega }_{r,p}^{}()`$.
We take $`d_1=p,d_2=r/p`$ in (3.2), and we see that $`\mathrm{\Omega }_{r,p}^{}()`$ is the image of the set $`\{0\}\times \mathrm{\Omega }_{r/p}^{}`$ under the isomorphism $`\xi `$ of (3.2). Therefore
$$\nu _{r,p}^{}=\nu _{r/p}^{}.$$
Using the above information, we find that
(6.3) $`_2`$ $`={\displaystyle \underset{d,e}{}}\lambda _{d,\mathrm{}_1}\lambda _{e,\mathrm{}_2}{\displaystyle \underset{p|[d,e]}{}}{\displaystyle \underset{a\mathrm{\Omega }_{[d,e],p}^{}}{}}{\displaystyle \underset{\begin{array}{c}N<n2N\\ na(mod[d,e])\end{array}}{}}\varpi \varpi (n)`$
$`=2N{\displaystyle \underset{d,e}{}}{\displaystyle \frac{\lambda _{d,\mathrm{}_1}\lambda _{e,\mathrm{}_2}}{\varphi ([d,e])}}{\displaystyle \underset{p|[d,e]}{}}\nu _{[d,e]/p}^{}{\displaystyle \frac{(\mathrm{log}p)\varphi (p)}{p}}+O\left((\mathrm{log}N)^{4k}_2\right),`$
where
$$_2=\stackrel{\mathrm{}}{}_{r<R^2}3^{\omega (r)}\underset{\begin{array}{c}p|r\\ p<R\end{array}}{}\nu _{r/p}^{}E^{}(N/p,r/p)(\mathrm{log}p).$$
Upon writing $`r=pm`$ and changing the order of summation, we find that
$$_23\underset{p<R}{}\mathrm{log}p\stackrel{\mathrm{}}{}_{m<R^2/p}3^{\omega (m)}\nu _m^{}E^{}(N/p,m).$$
By Lemma 2, the inner sum is $`(N/p)(\mathrm{log}N/p)^{4k2}(N/p)(\mathrm{log}N)^{4k2}`$. Summing over $`p`$, we get
$$_2N(\mathrm{log}N)^{4k1}.$$
Therefore
(6.4)
$$_2=2N\underset{d,e}{}\frac{\lambda _{d,\mathrm{}_1}\lambda _{e,\mathrm{}_2}}{\varphi ([d,e])}\underset{p|[d,e]}{}\frac{\nu _{[d,e]/p}^{}(\mathrm{log}p)\varphi (p)}{p}+O(N/\mathrm{log}N).$$
Now we turn our attention to $`_1`$. From our definitions and (1.10), we have
(6.5) $`_1=`$ $`{\displaystyle \underset{d,e}{}}\lambda _{d,\mathrm{}_1}\lambda _{e,\mathrm{}_2}{\displaystyle \underset{a\mathrm{\Omega }_{[d,e]}^{}}{}}{\displaystyle \underset{\begin{array}{c}Nn<2N\\ na(mod[d,e])\end{array}}{}}\varpi \varpi (n)`$
$`=`$ $`N{\displaystyle \underset{d,e}{}}{\displaystyle \frac{\lambda _{d,\mathrm{}_1}\lambda _{e,\mathrm{}_2}\nu _{[d,e]}^{}}{\varphi ([d,e])}}\left(\mathrm{log}N+C_02{\displaystyle \underset{p|[d,e]}{}}{\displaystyle \frac{\mathrm{log}p}{p}}\right)+O(_1),`$
where
$$_1=(\mathrm{log}R)^{4k}\stackrel{\mathrm{}}{}_{r<R^2}3^{\omega (r)}\nu _r^{}E_2^{}(N,r).$$
By Lemma 2, $`_1N/\mathrm{log}N`$.
Combining our estimates for $`_1`$ and $`_2`$, we find that
(6.6)
$$=N(\mathrm{log}N+C_0)S_12NS_2+2NS_3+O(N/\mathrm{log}N),$$
where
$`S_1=`$ $`{\displaystyle \underset{d,e}{}}{\displaystyle \frac{\lambda _{d,\mathrm{}_1}\lambda _{e,\mathrm{}_2}\nu _{[d,e]}^{}}{\varphi ([d,e])}},`$
$`S_2=`$ $`{\displaystyle \underset{d,e}{}}{\displaystyle \frac{\lambda _{d,\mathrm{}_1}\lambda _{e,\mathrm{}_2}\nu _{[d,e]}^{}}{\varphi ([d,e])}}{\displaystyle \underset{p|[d,e]}{}}{\displaystyle \frac{\mathrm{log}p}{p}},`$
$`S_3=`$ $`{\displaystyle \underset{d,e}{}}{\displaystyle \frac{\lambda _{d,\mathrm{}_1}\lambda _{e,\mathrm{}_2}}{\varphi ([d,e])}}{\displaystyle \underset{p|[d,e]}{}}\nu _{[d,e]/p}^{}{\displaystyle \frac{(\mathrm{log}p)\varphi (p)}{p}}.`$
We have already encountered the sum $`S_1`$; it is the same as the sum $`S`$ defined in (4.6). From (4.27), we see that
(6.7)
$$S_1=\left(\genfrac{}{}{0pt}{}{\mathrm{}_1+\mathrm{}_2+2}{\mathrm{}_1+1}\right)𝔖()\frac{(\mathrm{log}R)^{k+\mathrm{}_1+\mathrm{}_2+1}}{(k+\mathrm{}_1+\mathrm{}_2+1)!}+O\left(\beta ()𝔖()^2(\mathrm{log}R)^{k+\mathrm{}_1+\mathrm{}_2}\right).$$
Of the remaining two sums, $`S_3`$ is more important, so we concentrate on it first. We begin by interchanging the order of summation in $`S_3`$; this yields
(6.8)
$$S_3=\underset{p}{}\frac{\mathrm{log}p}{p}U(p),$$
where
(6.9)
$$U(p)=\underset{\begin{array}{c}d,e\\ p|[d,e]\end{array}}{}\frac{\lambda _{d,\mathrm{}_1}\lambda _{e,\mathrm{}_2}\nu _{[d,e]/p}^{}}{\varphi ([d,e]/p)}.$$
We decompose $`U(p)`$ as
(6.10)
$$U(p)=U_1(p)+U_2(p)+U_3(p),$$
where
$`U_1(p,\mathrm{}_1,\mathrm{}_2)`$ $`={\displaystyle \underset{\begin{array}{c}d,e\\ p|d,pe\end{array}}{}}{\displaystyle \frac{\lambda _{d,\mathrm{}_1}\lambda _{e,\mathrm{}_2}\nu _{[d,e]/p}^{}}{\varphi ([d,e]/p)}},`$
$`U_2(p,\mathrm{}_1,\mathrm{}_2)`$ $`={\displaystyle \underset{\begin{array}{c}d,e\\ pd,p|e\end{array}}{}}{\displaystyle \frac{\lambda _{d,\mathrm{}_1}\lambda _{e,\mathrm{}_2}\nu _{[d,e]/p}^{}}{\varphi ([d,e]/p)}},`$
$`U_3(p,\mathrm{}_1,\mathrm{}_2)`$ $`={\displaystyle \underset{\begin{array}{c}d,e\\ p|d,p|e\end{array}}{}}{\displaystyle \frac{\lambda _{d,\mathrm{}_1}\lambda _{e,\mathrm{}_2}\nu _{[d,e]/p}^{}}{\varphi ([d,e]/p)}}.`$
Going back to (6.8), we will write
(6.11)
$$S_3=S_{3,1}+S_{3,2}+S_{3,3},$$
where
$$S_{3,i}=\underset{p}{}\frac{\mathrm{log}p}{p}U_i(p,\mathrm{}_1,\mathrm{}_2).$$
We will ultimately see that each $`S_{3,i}`$ corresponds to one of the terms in the quantity $`T(k,\mathrm{}_1,\mathrm{}_2)`$ defined in the statement of Theorem 7. More precisely, we will show that when $`1i3`$,
$$S_{3,i}=T_i𝔖()\frac{(\mathrm{log}R)^{k+\mathrm{}_1+\mathrm{}_2+2}}{(k+\mathrm{}_1+\mathrm{}_2+2)!}\left\{1+O(\beta ()𝔖()/\mathrm{log}R)\right\},$$
where
$$T_1=\left(\genfrac{}{}{0pt}{}{\mathrm{}_1+\mathrm{}_2+3}{\mathrm{}_2+1}\right),T_2=\left(\genfrac{}{}{0pt}{}{\mathrm{}_1+\mathrm{}_2+3}{\mathrm{}_1+1}\right),T_3=\left(\genfrac{}{}{0pt}{}{\mathrm{}_1+\mathrm{}_2+2}{\mathrm{}_1+1}\right).$$
We note that $`U_2(p,\mathrm{}_1,\mathrm{}_2)`$ is the same as $`U_1(p,\mathrm{}_1,\mathrm{}_2)`$ except that the roles of $`\mathrm{}_1,\mathrm{}_2`$ have been reversed; i.e., $`U_2(p,\mathrm{}_1,\mathrm{}_2)=U_1(p,\mathrm{}_2,\mathrm{}_1)`$. Accordingly, we will concentrate on evaluating $`U_1(p,\mathrm{}_1,\mathrm{}_2)`$ and $`U_3(p,\mathrm{}_1,\mathrm{}_2)`$. For brevity, we will usually write these as $`U_1(p)`$ and $`U_3(p)`$.
The evaluations of $`U_1(p)`$ and $`U_3(p)`$ will require use of the quantity $`y_{r,\mathrm{}}^{}`$ defined in (4.11), as well as a new quantity $`z_{r,p,\mathrm{}}^{}`$. The latter is defined as
(6.12)
$$z_{r,p,\mathrm{}}^{}=\{\begin{array}{cc}\mu (pr)f_1^{}(r)\stackrel{}{}_d\frac{\lambda _{drp,\mathrm{}}}{f^{}(dr)}\hfill & \text{ if }r<R/p\text{ and }(r,A)=1\text{,}\hfill \\ 0\hfill & \text{ otherwise.}\hfill \end{array}$$
As in Section 4, we use $`_{}^{}{}_{}{}^{}`$ to denote that the sum is over values of the indices that are relatively prime to $`A`$. Note that $`z_{r,p,\mathrm{}}^{}=0`$ if $`(p,r)1`$. On the other hand, the condition $`p|A`$ (i.e., $`\nu _p^{}=0`$) does not imply that $`z_{r,p,\mathrm{}}^{}=0`$. However, one can easily show that if $`pA`$, then
(6.13)
$$z_{r,p,\mathrm{}}^{}=\left(\frac{p1}{p\nu _p}\right)y_{rp,\mathrm{}}^{}.$$
We now give three lemmas that we will use for the evaluation of $`S_1`$ and $`S_3`$.
###### Lemma 13.
If $`p<R`$, then
(6.14) $`U_1(p)`$ $`={\displaystyle \stackrel{}{}_{\begin{array}{c}r\\ (r,p)=1\end{array}}}{\displaystyle \frac{z_{r,p,\mathrm{}_1}^{}y_{r,\mathrm{}_2}^{}}{f_1^{}(r)}}{\displaystyle \frac{\nu _p^{}}{p1}}{\displaystyle \stackrel{}{}_{\begin{array}{c}r\\ (r,p)=1\end{array}}}{\displaystyle \frac{z_{r,p,\mathrm{}_1}^{}z_{r,p,\mathrm{}_2}^{}}{f_1^{}(r)}},\text{ and}`$
(6.15) $`U_3(p)`$ $`={\displaystyle \stackrel{}{}_{\begin{array}{c}r\\ (r,p)=1\end{array}}}{\displaystyle \frac{z_{r,p,\mathrm{}_1}^{}z_{r,p,\mathrm{}_2}^{}}{f_1^{}(r)}}.`$
###### Proof.
The sum $`U_1(p)`$ may be written as
(6.16) $`U_1(p)`$ $`={\displaystyle \stackrel{}{}_{\begin{array}{c}d,e\\ pe\end{array}}}{\displaystyle \frac{\lambda _{dp,\mathrm{}_1}\lambda _{e,\mathrm{}_2}}{\varphi ([d,e])}}\nu _{[d,e]}^{}={\displaystyle \stackrel{}{}_{\begin{array}{c}d,e\\ pe\end{array}}}{\displaystyle \frac{\lambda _{dp,\mathrm{}_1}\lambda _{e,\mathrm{}_2}}{f^{}([d,e])}}={\displaystyle \stackrel{}{}_{\begin{array}{c}d,e\\ pe\end{array}}}{\displaystyle \frac{\lambda _{dp,\mathrm{}_1}\lambda _{e,\mathrm{}_2}}{f^{}(d)f^{}(e)}}{\displaystyle \underset{\begin{array}{c}r|d\\ r|e\end{array}}{}}f_1^{}(r)`$
$`={\displaystyle \stackrel{}{}_{\begin{array}{c}r\\ (r,p)=1\end{array}}}f_1^{}(r)\left({\displaystyle \underset{d}{}}{\displaystyle \frac{\lambda _{drp,\mathrm{}_1}}{f^{}(dr)}}\right)\left({\displaystyle \underset{\begin{array}{c}e\\ pe\end{array}}{}}{\displaystyle \frac{\lambda _{er,\mathrm{}_2}}{f^{}(er)}}\right).`$
In the last expression, the first sum in parentheses is $`\mu (pr)z_{r,p,\mathrm{}_1}^{}/f_1^{}(r)`$. The innermost sum is
$$\stackrel{}{}_e\frac{\lambda _{er,\mathrm{}_2}}{f^{}(er)}\stackrel{}{}_{\begin{array}{c}e\\ p|e\end{array}}\frac{\lambda _{er,\mathrm{}_2}}{f^{}(er)}=\frac{\mu (r)y_{r,\mathrm{}_2}^{}}{f_1^{}(r)}\stackrel{}{}_{\begin{array}{c}e\\ p|e\end{array}}\frac{\lambda _{er,\mathrm{}_2}}{f^{}(er)}.$$
We claim that
(6.17)
$$\stackrel{}{}_{\begin{array}{c}e\\ p|e\end{array}}\frac{\lambda _{er,\mathrm{}_2}}{f^{}(er)}=\frac{\nu _p^{}\mu (pr)z_{r,p,\mathrm{}_2}^{}}{(p1)f_1^{}(r)}.$$
If $`\nu _p^{}=0`$, then both sides of (6.17) are $`0`$. If $`\nu _p^{}0`$, then
$$\stackrel{}{}_{\begin{array}{c}e\\ p|e\end{array}}\frac{\lambda _{er,\mathrm{}_2}}{f^{}(er)}=\stackrel{}{}_e\frac{\lambda _{epr,\mathrm{}_2}}{f^{}(epr)}=\frac{\mu (pr)z_{r,p,\mathrm{}_2}^{}}{f_1^{}(r)f^{}(p)},$$
and (6.17) follows again.
Going back to (6.16), we find that
$$U_1(p)=\stackrel{}{}_{\begin{array}{c}r\\ (r,p)=1\end{array}}f_1^{}(r)\left(\frac{\mu (rp)z_{r,p,\mathrm{}_1}^{}}{f_1^{}(r)}\right)\left(\frac{\mu (r)y_{r,\mathrm{}_2}^{}}{f_1^{}(r)}\frac{\mu (rp)z_{r,p,\mathrm{}_2}^{}\nu _p^{}}{f_1^{}(r)(p1)}\right),$$
and (6.14) follows.
For $`U_3(p)`$, observe that
$`U_3(p)=`$ $`{\displaystyle \stackrel{}{}_{d,e}}{\displaystyle \frac{\lambda _{dp,\mathrm{}_1}\lambda _{ep,\mathrm{}_2}}{f^{}([d,e])}}={\displaystyle \stackrel{}{}_{d,e}}{\displaystyle \frac{\lambda _{dp,\mathrm{}_1}\lambda _{ep,\mathrm{}_2}}{f^{}(d)f^{}(e)}}{\displaystyle \underset{\begin{array}{c}r|d\\ r|e\end{array}}{}}f_1^{}(r)`$
$`=`$ $`{\displaystyle \stackrel{}{}_{\begin{array}{c}r\\ (r,p)=1\end{array}}}f_1^{}(r)\left({\displaystyle \stackrel{}{}_d}{\displaystyle \frac{\lambda _{drp,\mathrm{}_1}}{f^{}(dr)}}\right)\left({\displaystyle \stackrel{}{}_e}{\displaystyle \frac{\lambda _{erp,\mathrm{}_1}}{f^{}(er)}}\right)`$
$`=`$ $`{\displaystyle \stackrel{}{}_{\begin{array}{c}r\\ (r,p)=1\end{array}}}{\displaystyle \frac{z_{r,p,\mathrm{}_1}^{}z_{r,p,\mathrm{}_2}^{}}{f_1^{}(r)}},`$
and this yields (6.15). ∎
###### Lemma 14.
If $`r<R/p`$ and $`(r,A)=1`$, then
(6.18) $`z_{r,p,\mathrm{}}^{}=\mu ^2(rp)`$ $`{\displaystyle \frac{𝔖()}{(\mathrm{}+1)!}}\left({\displaystyle \frac{p1}{p\nu _p}}\right)(\mathrm{log}R/rp)^{\mathrm{}+1}`$
$`+O\left(\mu ^2(rp){\displaystyle \frac{\rho (rp)rp}{\varphi (rp)}}𝔖()(\mathrm{log}2R/rp)^{\mathrm{}}\right).`$
We remark that the error term could be simplified; it is obvious that
$$\frac{\rho (rp)rp}{\varphi (rp)}\frac{\rho (r)r}{\varphi (r)}.$$
However, we prefer to write it as above to emphasize the connection between $`y_{r,\mathrm{}}^{}`$ and $`z_{r,p,\mathrm{}}^{}`$. In fact, this lemma follows immediately from (6.13) and (4.17) when $`\nu _p^{}0`$. However, the following argument works whether or not $`\nu _p^{}=0`$.
###### Proof.
The result is trivial is $`rp`$ is not squarefree, becuase both sides of (6.18) are 0 in this case. For the rest of this proof, we assume that $`rp`$ is squarefree. Note that this assumption implies that $`(r,p)=1`$.
We start by observing that
$`{\displaystyle \frac{\mu (rp)z_{r,p,\mathrm{}}^{}}{f_1^{}(r)}}`$ $`={\displaystyle \stackrel{}{}_d}{\displaystyle \frac{\lambda _{drp,\mathrm{}}}{f^{}(dr)}}={\displaystyle \stackrel{}{}_d}{\displaystyle \frac{\mu (drp)}{f^{}(dr)}}f(drp){\displaystyle \underset{t}{}}{\displaystyle \frac{y_{drpt,\mathrm{}}}{f_1(drpt)}}`$
$`={\displaystyle \frac{\mu (rp)f(rp)}{f^{}(r)f_1(rp)}}{\displaystyle \stackrel{}{}_d}{\displaystyle \frac{\mu (d)f(d)}{f^{}(d)}}{\displaystyle \underset{t}{}}{\displaystyle \frac{y_{rpdt,\mathrm{}}}{f_1(dt)}}`$
$`={\displaystyle \frac{\mu (rp)f(rp)}{f^{}(r)f_1(rp)}}{\displaystyle \underset{\begin{array}{c}m\\ (m,rp)=1\end{array}}{}}{\displaystyle \frac{y_{rpm,\mathrm{}}}{f_1(m)}}{\displaystyle \stackrel{}{}_{d|m}}{\displaystyle \frac{\mu (d)f(d)}{f^{}(d)}}`$
$`={\displaystyle \frac{\mu (rp)f(rp)}{f^{}(r)f_1(rp)}}{\displaystyle \underset{\begin{array}{c}m\\ (m,rp)=1\end{array}}{}}{\displaystyle \frac{y_{rpm,\mathrm{}}}{\varphi (m)}}.`$
In the last line, we have used the relation (4.12). If we also use (4.13), we find that
(6.19) $`z_{r,p,\mathrm{}}^{}`$ $`={\displaystyle \frac{f(r)f_1^{}(r)f(p)}{f_1(r)f^{}(r)f_1(p)}}{\displaystyle \underset{\begin{array}{c}m<R/rp\\ (m,rp)=1\end{array}}{}}{\displaystyle \frac{y_{rpm,\mathrm{}}}{\varphi (m)}}`$
$`={\displaystyle \frac{𝔖()}{\mathrm{}!}}{\displaystyle \frac{rp}{\varphi (rp)}}{\displaystyle \frac{(p1)}{(p\nu _p)}}{\displaystyle \underset{\begin{array}{c}m<R/rp\\ (m,rp)=1\end{array}}{}}{\displaystyle \frac{\mu ^2(m)}{\varphi (m)}}(\mathrm{log}R/rpm)^{\mathrm{}}.`$
We then use (4.16) to complete the proof. ∎
###### Lemma 15.
If $`a,b`$ are non-negative integers, then
$`{\displaystyle \underset{p<R}{}}{\displaystyle \frac{(\mathrm{log}p)^{a+1}(\mathrm{log}R/p)^b}{p}}=`$ $`{\displaystyle \frac{a!b!}{(a+b+1)!}}(\mathrm{log}R)^{a+b+1}+O_{a,b}((\mathrm{log}R)^{a+b}),\text{ and}`$
$`{\displaystyle \underset{p<R}{}}{\displaystyle \frac{(\mathrm{log}p)^{a+1}(\mathrm{log}R/p)^b}{p^2}}`$ $`_a(\mathrm{log}R)^b.`$
###### Proof.
Let $`E(u)`$ be defined by the relation
$$\underset{pu}{}\frac{\mathrm{log}p}{p}=\mathrm{log}u+E(u).$$
It is well-known that $`E(u)1`$. The first sum in the lemma is
$`{\displaystyle \underset{p<R}{}}`$ $`{\displaystyle \frac{(\mathrm{log}p)^{a+1}(\mathrm{log}R/p)^b}{p}}=`$
$`{\displaystyle _1^R}(\mathrm{log}u)^a(\mathrm{log}R/u)^b{\displaystyle \frac{du}{u}}+{\displaystyle _1^R}(\mathrm{log}u)^a(\mathrm{log}R/u)^b𝑑E(u).`$
By Lemma 3, the first integral is
$$\frac{a!b!}{(a+b+1)!}(\mathrm{log}R)^{a+b+1}.$$
Using integration by parts, we see that the second integral is
$$_1^RE(u)\frac{d}{du}\left\{(\mathrm{log}u)^a(\mathrm{log}R/u)^b\right\}𝑑u_{a,b}(\mathrm{log}R)^{a+b}.$$
This proves the first statement. The second statement is easier; we simply note that
$$\underset{p<R}{}\frac{(\mathrm{log}p)^{a+1}(\mathrm{log}R/p)^b}{p^2}(\mathrm{log}R)^b\underset{p}{}\frac{(\mathrm{log}p)^{a+1}}{p^2}_a(\mathrm{log}R)^b.$$
Evaluation of $`S_{3,3}`$. From Lemmas 13 and 14, we see that
(6.20)
$$U_3(p)=\frac{𝔖()^2}{(\mathrm{}_1+1)!(\mathrm{}_2+1)!}\left(\frac{p1}{p\nu _p}\right)^2V_3(p)+O\left(𝔖()^2(\mathrm{log}R)^{\mathrm{}_1+\mathrm{}_2+1}W(p)\right),$$
where
(6.21) $`V_3(p)`$ $`={\displaystyle \stackrel{}{}_{\begin{array}{c}r<R/p\\ (r,p)=1\end{array}}}{\displaystyle \frac{\mu ^2(r)}{f_1^{}(r)}}(\mathrm{log}R/rp)^{\mathrm{}_1+\mathrm{}_2+2},\text{ and}`$
(6.22) $`W(p)`$ $`={\displaystyle \stackrel{}{}_{r<R/p}}{\displaystyle \frac{\mu ^2(r)\rho (r)r}{f_1^{}(r)\varphi (r)}}.`$
$`W(p)`$ is majorized by the sum $`W`$ defined in (4.20), and, using (4.26), we see that
(6.23)
$$W(p)(\mathrm{log}R)^{k1}.$$
From Lemma 10, we see that
(6.24)
$$V_3(p)=\frac{(\mathrm{}_1+\mathrm{}_2+2)!}{𝔖()}\left(\frac{p\nu _p}{p1}\right)\frac{(\mathrm{log}R/p)^{k+\mathrm{}_1+\mathrm{}_2+1}}{(k+\mathrm{}_1+\mathrm{}_2+1)!}+O\left(\beta ()(\mathrm{log}R)^{k+\mathrm{}_1+\mathrm{}_2}\right).$$
We combine the above estimates for $`V_3(p)`$ and $`W(p)`$ with (6.20) to get
(6.25) $`U_3(p)=\left({\displaystyle \genfrac{}{}{0pt}{}{\mathrm{}_1+\mathrm{}_2+2}{\mathrm{}_1+1}}\right)𝔖()`$ $`\left({\displaystyle \frac{p1}{p\nu _p}}\right){\displaystyle \frac{(\mathrm{log}R/p)^{k+\mathrm{}_1+\mathrm{}_2+1}}{(k+\mathrm{}_1+\mathrm{}_2+1)!}}`$
$`+O\left(\beta ()𝔖()^2(\mathrm{log}R)^{k+\mathrm{}_1+\mathrm{}_2}\right).`$
We can now finish our estimation of $`S_{3,3}`$. From our definition and from (6.25), we get
$`S_{3,3}=`$ $`{\displaystyle \underset{p<R}{}}{\displaystyle \frac{\mathrm{log}p}{p}}U_3(p)`$
$`=`$ $`\left({\displaystyle \genfrac{}{}{0pt}{}{\mathrm{}_1+\mathrm{}_2+2}{\mathrm{}_1+1}}\right){\displaystyle \frac{𝔖()}{(k+\mathrm{}_1+\mathrm{}_2+1)!}}{\displaystyle \underset{p<R}{}}\left({\displaystyle \frac{\mathrm{log}p}{p}}\right)\left({\displaystyle \frac{p1}{p\nu _p}}\right)(\mathrm{log}R/p)^{k+\mathrm{}_1+\mathrm{}_2+1}`$
$`+O\left(\beta ()𝔖()^2(\mathrm{log}R)^{k+\mathrm{}_1+\mathrm{}_2}{\displaystyle \underset{p<R}{}}{\displaystyle \frac{\mathrm{log}p}{p}}\right).`$
Now $`(p1)/(p\nu _p)=1+O(1/p)`$, so we may use Lemma 15 to get
(6.26) $`S_{3,3}=`$ $`\left({\displaystyle \genfrac{}{}{0pt}{}{\mathrm{}_1+\mathrm{}_2+2}{\mathrm{}_1+1}}\right)𝔖(){\displaystyle \frac{(\mathrm{log}R)^{k+\mathrm{}_1+\mathrm{}_2+2}}{(k+\mathrm{}_1+\mathrm{}_2+2)!}}`$
$`+O\left(\beta ()𝔖()^2(\mathrm{log}R)^{k+\mathrm{}_1+\mathrm{}_2+1}\right).`$
Evaluation of $`S_{3,1}`$. The evaluation of $`S_{3,1}`$ proceeds similarly to the evaluation of $`S_{3,3}`$, but it is somewhat more involved. We start by defining
(6.27)
$$U_4(p)=\stackrel{}{}_{\begin{array}{c}r\\ (r,p)=1\end{array}}\frac{y_{r,\mathrm{}_2}^{}z_{r,p,\mathrm{}_1}^{}}{f_1^{}(r)},$$
and
(6.28)
$$S_4=\underset{p<R}{}\frac{\mathrm{log}p}{p}U_4(p).$$
Then (6.14) may be rewritten as
(6.29)
$$U_1(p)=U_4(p)\frac{\nu _p^{}}{p1}U_3(p),$$
and we may also write
(6.30)
$$S_{3,1}=S_4\underset{p<R}{}\frac{(\mathrm{log}p)\nu _p^{}}{p(p1)}U_3(p).$$
From (6.25), we see that
(6.31) $`{\displaystyle \underset{p<R}{}}{\displaystyle \frac{(\mathrm{log}p)\nu _p^{}}{p(p1)}}U_3(p)`$ $`\beta ()𝔖()^2(\mathrm{log}R)^{k+\mathrm{}_1+\mathrm{}_2+1}{\displaystyle \underset{p}{}}{\displaystyle \frac{(\mathrm{log}p)\nu _p^{}}{p^2}}`$
$`\beta ()𝔖()^2(\mathrm{log}R)^{k+\mathrm{}_1+\mathrm{}_2+1}.`$
Now we concentrate on $`U_4(p)`$ and $`S_4`$. From (4.17) and Lemma 14, we see that
(6.32)
$$U_4(p)=\frac{𝔖()^2}{(\mathrm{}_1+1)!(\mathrm{}_2+1)!}\left(\frac{p1}{p\nu _p}\right)V_4(p)+O\left(𝔖()^2(\mathrm{log}R)^{\mathrm{}_1+\mathrm{}_2+1}W(p)\right),$$
where $`W(p)`$ was defined in (6.22) and
(6.33)
$$V_4(p)=\stackrel{}{}_{\begin{array}{c}r<R/p\\ (r,p)=1\end{array}}\frac{\mu ^2(r)}{f_1^{}(r)}(\mathrm{log}R/r)^{\mathrm{}_2+1}(\mathrm{log}R/rp)^{\mathrm{}_1+1}.$$
We write $`\mathrm{log}R/r=\mathrm{log}p+\mathrm{log}R/rp`$ and use the binomial theorem to get
$$V_4(p)=\underset{j=0}{\overset{\mathrm{}_2+1}{}}\left(\genfrac{}{}{0pt}{}{\mathrm{}_2+1}{j}\right)(\mathrm{log}p)^j\stackrel{}{}_{\begin{array}{c}r<R/p\\ (r,p)=1\end{array}}\frac{\mu ^2(r)}{f_1^{}(r)}(\mathrm{log}R/rp)^{\mathrm{}_1+\mathrm{}_2+2j}.$$
We apply Lemma 10 to the inner sum, and we get
(6.34) $`V_4(p)=`$
$`{\displaystyle \frac{1}{𝔖()}}\left({\displaystyle \frac{p\nu _p}{p1}}\right){\displaystyle \underset{j=0}{\overset{\mathrm{}_2+1}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{\mathrm{}_2+1}{j}}\right){\displaystyle \frac{(\mathrm{}_1+\mathrm{}_2+2j)!}{(k+\mathrm{}_1+\mathrm{}_2+1j)!}}(\mathrm{log}p)^j(\mathrm{log}R/p)^{k+\mathrm{}_1+\mathrm{}_2+1j}`$
$`+O(\beta ()(\mathrm{log}2R)^{k+\mathrm{}_1+\mathrm{}_2}).`$
Using this together with (6.32) and (6.23) gives
(6.35)
$$U_4(p)=\frac{𝔖()}{(\mathrm{}_1+1)!(\mathrm{}_2+1)!}U_5(p)+O(\beta ()𝔖()^2(\mathrm{log}R)^{k+\mathrm{}_1+\mathrm{}_2}),$$
where
(6.36)
$$U_5(p)=\underset{j=0}{\overset{\mathrm{}_2+1}{}}\left(\genfrac{}{}{0pt}{}{\mathrm{}_2+1}{j}\right)\frac{(\mathrm{}_1+\mathrm{}_2+2j)!}{(k+\mathrm{}_1+\mathrm{}_2+1j)!}(\mathrm{log}p)^j(\mathrm{log}R/p)^{k+\mathrm{}_1+\mathrm{}_2+1j}.$$
For future reference, we note the crude estimate
(6.37)
$$U_1(p)\beta ()𝔖()^2(\mathrm{log}R)^{k+\mathrm{}_1+\mathrm{}_2+1}$$
that is implicit in the combination of (6.29), (6.35), (6.36), and (6.25).
Using (6.28) and (6.35), we see that
(6.38)
$$S_4=\frac{𝔖()}{(\mathrm{}_1+1)!(\mathrm{}_2+1)!}\underset{p<R}{}\frac{\mathrm{log}p}{p}U_5(p)+O(\beta ()𝔖()^2(\mathrm{log}R)^{k+\mathrm{}_1+\mathrm{}_2+1}).$$
We apply Lemma 15 to get
$`{\displaystyle \underset{p<R}{}}`$ $`\left({\displaystyle \frac{\mathrm{log}p}{p}}\right)(\mathrm{log}p)^j(\mathrm{log}R/p)^{k+\mathrm{}_1+\mathrm{}_2+1j}`$
$`=`$ $`{\displaystyle \frac{j!(k+\mathrm{}_1+\mathrm{}_2+1j)!}{(k+\mathrm{}_1+\mathrm{}_2+2)!}}(\mathrm{log}R)^{k+\mathrm{}_1+\mathrm{}_2+2}+O((\mathrm{log}R)^{k+\mathrm{}_1+\mathrm{}_1+1}).`$
Using this in (6.38) gives
(6.39) $`S_4=𝔖()`$ $`{\displaystyle \frac{(\mathrm{log}R)^{k+\mathrm{}_1+\mathrm{}_2+2}}{(k+\mathrm{}_1+\mathrm{}_2+2)!}}{\displaystyle \underset{j=0}{\overset{\mathrm{}_2+1}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{\mathrm{}_1+\mathrm{}_2+2j}{\mathrm{}_2+1j}}\right)`$
$`+O(\beta ()𝔖()^2(\mathrm{log}R)^{k+\mathrm{}_1+\mathrm{}_2+1}).`$
To treat the sum of binomial coefficients in the above, we make a change of variables $`j=\mathrm{}_2+1i`$. The sum then becomes
(6.40)
$$\underset{i=0}{\overset{\mathrm{}_2+1}{}}\left(\genfrac{}{}{0pt}{}{\mathrm{}_1+1+i}{i}\right)=\underset{i=0}{\overset{\mathrm{}_2+1}{}}\left\{\left(\genfrac{}{}{0pt}{}{\mathrm{}_1+2+i}{i}\right)\left(\genfrac{}{}{0pt}{}{\mathrm{}_1+1+i}{i1}\right)\right\},$$
provided we make the usual convention that
$$\left(\genfrac{}{}{0pt}{}{\mathrm{}_1+1}{1}\right)=0.$$
The sum on the right-hand side of (6.40) is telescoping, so
$$\underset{i=0}{\overset{\mathrm{}_2+1}{}}\left(\genfrac{}{}{0pt}{}{\mathrm{}_1+1+i}{i}\right)=\left(\genfrac{}{}{0pt}{}{\mathrm{}_1+\mathrm{}_2+3}{\mathrm{}_2+1}\right).$$
Putting this information into (6.39) gives our final estimate for $`S_4`$; i.e.,
(6.41) $`S_4=\left({\displaystyle \genfrac{}{}{0pt}{}{\mathrm{}_1+\mathrm{}_2+3}{\mathrm{}_2+1}}\right)`$ $`𝔖(){\displaystyle \frac{(\mathrm{log}R)^{k+\mathrm{}_1+\mathrm{}_2+2}}{(k+\mathrm{}_1+\mathrm{}_2+2)!}}`$
$`+O(\beta ()𝔖()^2(\mathrm{log}R)^{k+\mathrm{}_1+\mathrm{}_2+1}).`$
From this, together with (6.30) and (6.31), we get
(6.42) $`S_{3,1}=\left({\displaystyle \genfrac{}{}{0pt}{}{\mathrm{}_1+\mathrm{}_2+3}{\mathrm{}_2+1}}\right)`$ $`𝔖(){\displaystyle \frac{(\mathrm{log}R)^{k+\mathrm{}_1+\mathrm{}_2+2}}{(k+\mathrm{}_1+\mathrm{}_2+2)!}}`$
$`+O(\beta ()𝔖()^2(\mathrm{log}R)^{k+\mathrm{}_1+\mathrm{}_2+1}).`$
As we noted earlier, $`S_{3,2}`$ is the same as $`S_{3,1}`$ with the roles of $`\mathrm{}_1`$ and $`\mathrm{}_2`$ reversed. Therefore
(6.43) $`S_{3,2}=\left({\displaystyle \genfrac{}{}{0pt}{}{\mathrm{}_1+\mathrm{}_2+3}{\mathrm{}_1+1}}\right)`$ $`𝔖(){\displaystyle \frac{(\mathrm{log}R)^{k+\mathrm{}_1+\mathrm{}_2+2}}{(k+\mathrm{}_1+\mathrm{}_2+2)!}}`$
$`+O(\beta ()𝔖()^2(\mathrm{log}R)^{k+\mathrm{}_1+\mathrm{}_2+1}).`$
Combining (6.42),(6.43), and (6.26) gives
(6.44) $`S_3=T(k,\mathrm{}_1,\mathrm{}_2)`$ $`𝔖(){\displaystyle \frac{(\mathrm{log}R)^{k+\mathrm{}_1+\mathrm{}_2+2}}{(k+\mathrm{}_1+\mathrm{}_2+2)!}}`$
$`+O(\beta ()𝔖()^2(\mathrm{log}R)^{k+\mathrm{}_1+\mathrm{}_2+1}),`$
where $`T(k,\mathrm{}_1,\mathrm{}_2)`$ is as defined in Theorem 7.
Finally, we will quickly dispatch $`S_2`$. We rewrite this sum as
$$S_2=\stackrel{}{}_{d,e}\frac{\lambda _{d,\mathrm{}_1}\lambda _{e,\mathrm{}_2}}{f^{}([d,e])}\underset{p|[d,e]}{}\frac{\mathrm{log}p}{p}=\stackrel{}{}_p\frac{\mathrm{log}p}{pf^{}(p)}U(p),$$
where $`U(p)`$ was defined in (6.9). We employ the crude estimate
$$U(p)𝔖()^2\beta ()(\mathrm{log}R)^{k+\mathrm{}_1+\mathrm{}_2+1}.$$
This is easily seen by combining (6.10), (6.37), (6.25), and using the symmetry between $`U_1(p)`$ and $`U_2(p)`$. The sum
$$\stackrel{}{}_{pR}\frac{\mathrm{log}p}{pf^{}(p)}$$
is $`1`$. Combining the above gives the bound
(6.45)
$$S_2𝔖()^2\beta ()(\mathrm{log}R)^{k+\mathrm{}_1+\mathrm{}_2+1}.$$
The proof of Theorem 7 is completed by combining (6.6) together with the final estimates for $`S_1,S_2,S_3`$, which are (6.7), (6.45), and (6.44) respectively.
## 7. Proofs of Theorems 1 through 4
Let $`=\{h_1,h_2,\mathrm{},h_k\}`$ be an arbitrary admissible $`k`$-tuple. Without loss of generality, we may specify that
$$h_1<h_2<\mathrm{}<h_k.$$
It is also useful to assume that
(7.1)
$$h_k\mathrm{log}N.$$
With this hypothesis, we see from Lemma 6 that the error terms in Theorems 5, 6, 7 satisfy
$$\beta ()𝔖()/\mathrm{log}N(\mathrm{log}\mathrm{log}\mathrm{log}N)^{b_k+1}/\mathrm{log}N(\mathrm{log}\mathrm{log}N)/\mathrm{log}N.$$
Consider the sum
(7.2)
$$𝒮_1:=\underset{N<n2N}{}\left\{\underset{h}{}\varpi (n+h)(\mathrm{log}3N)\right\}\left(\underset{\mathrm{}=0}{\overset{L}{}}b_{\mathrm{}}(\mathrm{log}R)^{\mathrm{}}\mathrm{\Lambda }_R(n;,\mathrm{})\right)^2.$$
For a given $`n`$, the sum inside the brackets is non-positive unless there are at least two distinct values, $`h_i,h_j`$ such that $`n+h_i,n+h_j`$ are primes. Consequently, if we can show that the sum in (7.2) is $`N𝔖()(\mathrm{log}R)^{k+1}`$, then we can conclude that $`lim\; inf_n\mathrm{}(p_{n+1}p_n)h_kh_1`$.
Expanding the square in (7.2), we see that
$$𝒮_1=\underset{0\mathrm{}_1,\mathrm{}_2L}{}b_\mathrm{}_1b_\mathrm{}_2(\mathrm{log}R)^{\mathrm{}_1\mathrm{}_2}_1(\mathrm{}_1,\mathrm{}_2),$$
where
$$_1(\mathrm{}_1,\mathrm{}_2)=\underset{Nn<2N}{}\left\{\underset{h}{}\varpi (n+h)(\mathrm{log}3N)\right\}\mathrm{\Lambda }_R(n;,\mathrm{}_1)\mathrm{\Lambda }_R(n;,\mathrm{}_2).$$
We assume Hypothesis $`BV(\theta )`$, and we use Theorems 5 and 6 with $`R=N^{(\theta ϵ)/2}`$ to get
$`_1(\mathrm{}_1,\mathrm{}_2)`$ $`\left({\displaystyle \genfrac{}{}{0pt}{}{\mathrm{}_1+\mathrm{}_2+2}{\mathrm{}_1+1}}\right)N𝔖()k{\displaystyle \frac{(\mathrm{log}R)^{k+\mathrm{}_1+\mathrm{}_2+1}}{(k+\mathrm{}_1+\mathrm{}_2+1)!}}`$
$`\left({\displaystyle \genfrac{}{}{0pt}{}{\mathrm{}_1+\mathrm{}_2}{\mathrm{}_1}}\right)N𝔖(){\displaystyle \frac{(\mathrm{log}R)^{k+\mathrm{}_1+\mathrm{}_2}\mathrm{log}N}{(k+\mathrm{}_1+\mathrm{}_2)!}}`$
$`N𝔖()(\mathrm{log}R)^{k+\mathrm{}_1+\mathrm{}_2}(\mathrm{log}N)(m(k,\mathrm{}_1,\mathrm{}_2,\theta )ϵ^{})`$
where
(7.3)
$$m(k,\mathrm{}_1,\mathrm{}_2,\theta )=\left(\genfrac{}{}{0pt}{}{\mathrm{}_1+\mathrm{}_2}{\mathrm{}_1}\right)\frac{1}{(k+\mathrm{}_1+\mathrm{}_2)!}\left(\frac{k(\mathrm{}_1+\mathrm{}_2+1)(\mathrm{}_1+\mathrm{}_2+2)}{(k+\mathrm{}_1+\mathrm{}_2+1)(\mathrm{}_1+1)(\mathrm{}_2+1)}\frac{\theta }{2}1\right),$$
and $`ϵ^{}=ϵ^{}(k,\mathrm{}_1,\mathrm{}_2,ϵ)`$ goes to 0 as $`ϵ`$ goes to $`0`$.
Define $`𝐛=(b_0,b_1,\mathrm{},b_L)`$. Then (we suppress the $`ϵ^{}`$ term)
(7.4) $`𝒮_1^{}(N,,\theta ,𝐛)`$ $`:={\displaystyle \frac{𝒮_1}{N𝔖()(\mathrm{log}R)^k\mathrm{log}N}}`$
$`{\displaystyle \underset{0\mathrm{}_1,\mathrm{}_2L}{}}b_\mathrm{}_1b_\mathrm{}_2m(k,\mathrm{}_1,\mathrm{}_2,\theta )`$
$`=𝐛^T\mathrm{𝐌𝐛},`$
where $`𝐌=𝐌(k,\theta )`$ is the matrix
$$𝐌=\left[m(k,i,j,\theta )\right]_{0i,jL}.$$
Our goal is to pick $`𝐛`$ to make $`𝒮_1^{}>0`$ for a given $`\theta `$ and minimal $`k`$. This is easily determined by picking $`𝐛`$ to be an eigenvector of the matrix $`𝐌`$ with eigenvalue $`\lambda `$, in which case
$$𝒮_1^{}𝐛^T\lambda 𝐛=\lambda \underset{i=0}{\overset{L}{}}b_i^2.$$
This will be positive provided $`\lambda `$ is positive. We conclude that $`𝒮_1^{}>0`$ if $`𝐌`$ has a positive eigenvalue and $`𝐛`$ is chosen to be the corresponding eigenvector.
With $`k=6`$ and $`L=1`$, we find that
$$𝐌=\frac{1}{8!}\left[\begin{array}{cc}48\theta 56& 9\theta 8\\ 9\theta 8& 2\theta 2\end{array}\right].$$
The determinant of $`8!𝐌`$ is $`15\theta ^264\theta +48`$, which is negative if $`4(8\sqrt{19})/15<\theta 1`$. Since the determinant is the product of the eigenvalues, we conclude that $`𝐌`$ has a positive eigenvalue for $`\theta `$ in this range. Consequently, if $``$ is an admissible $`6`$-tuple, then there are infinitely many $`n`$ such that at least two of the numbers $`n+h_1,\mathrm{},n+h_6`$ are prime. We complete the proof of the second part of Theorem 2 by taking
$$=\{7,11,13,17,19,23\}.$$
$``$ is admissible because for $`p5`$, none of the elements in $``$ are divisible by $`p`$, and for $`p7`$, there are not enough elements to cover all of the residue classes mod $`p`$.
To prove the first part of Theorem 2, we again use (7.4); however, we use the trivial choice $`b_{\mathrm{}}=1`$ for some specific $`\mathrm{}`$, and $`b_i=0`$ for all other $`i`$. Then
$$𝒮_1^{}m(k,\mathrm{},\mathrm{},\theta )=\left(\genfrac{}{}{0pt}{}{2\mathrm{}}{\mathrm{}}\right)\frac{1}{(k+2\mathrm{})!}\left(\frac{2k(2\mathrm{}+1)}{(k+2\mathrm{}+1)(\mathrm{}+1)}\frac{\theta }{2}1\right)ϵ^{}.$$
The above is positive if
$$\theta >\left(\frac{1}{2}+\frac{1}{4\mathrm{}+2}\right)\left(1+\frac{2\mathrm{}+1}{k}\right).$$
The right-hand side approaches $`1/2`$ if $`\mathrm{},k\mathrm{}`$ with $`\mathrm{}=o(k)`$.
The above argument just fails when $`\theta =1/2`$. To remedy this, we modify (7.2) by taking $`h`$ to be a parameter to be chosen later, with $`h\mathrm{log}N`$. We then sum over all admissible size $`k`$ subsets $``$ of $`\{1,\mathrm{},h\}`$. Specifically, we take
(7.5)
$$\stackrel{~}{𝒮}_1=\underset{\begin{array}{c}\{1,\mathrm{},h\}\\ ||=k\\ \text{ admissible}\end{array}}{}\underset{N<n2N}{}\left\{\underset{1h_0h}{}\varpi (n+h_0)(\mathrm{log}3N)\right\}\mathrm{\Lambda }_R^2(n;,\mathrm{}).$$
We apply Theorems 5 and 6 to the sum $`\stackrel{~}{𝒮}_1`$ for those terms when $``$ and $`\{h_0\}`$ are both admissible. There may be terms with $``$ admissible but $`\{h_0\}`$ not admissible; for these terms we apply the trivial bound
$$\underset{N<n2N}{}\underset{1h_0h}{}\varpi (n+h_0)\mathrm{\Lambda }_R(n;,\mathrm{})^20.$$
We find that
(7.6) $`\stackrel{~}{𝒮}_1`$ $`\left({\displaystyle \genfrac{}{}{0pt}{}{2\mathrm{}+2}{\mathrm{}+1}}\right){\displaystyle \frac{N(\mathrm{log}R)^{k+2\mathrm{}+1}}{(k+2\mathrm{}+1)!}}{\displaystyle \underset{1h_0h}{}}{\displaystyle \underset{\begin{array}{c}\{1,\mathrm{},h\}\\ ||=k,h_0\end{array}}{}}𝔖()`$
$`+\left({\displaystyle \genfrac{}{}{0pt}{}{2\mathrm{}}{\mathrm{}}}\right){\displaystyle \frac{N(\mathrm{log}R)^{k+2\mathrm{}}}{(k+2\mathrm{})!}}{\displaystyle \underset{1h_0h}{}}{\displaystyle \underset{\begin{array}{c}\{1,\mathrm{},h\}\\ ||=k,h_0\end{array}}{}}𝔖\left(\{h_0\}\right)`$
$`\left({\displaystyle \genfrac{}{}{0pt}{}{2\mathrm{}}{\mathrm{}}}\right){\displaystyle \frac{N(\mathrm{log}N)(\mathrm{log}R)^{k+2\mathrm{}}}{(k+2\mathrm{})!}}{\displaystyle \underset{\begin{array}{c}\{1,\mathrm{},h\}\\ ||=k\end{array}}{}}𝔖().`$
We have dropped the condition that $``$ is admissible in the above sums; we may do so because $`𝔖()=0`$ when $``$ is not admissible.
Now we observe that
$$\underset{1h_0h}{}\underset{\begin{array}{c}\{1,\mathrm{},h\}\\ ||=k,h_0\end{array}}{}𝔖()=k\underset{\begin{array}{c}\{1,\mathrm{},h\}\\ ||=k\end{array}}{}𝔖()\frac{kh^k}{k!}.$$
In the above, equality occurs from noting that every relevant set $``$ occurs $`k`$ times in the initial sum, and the asymptotic relation is a theorem of Gallagher . We also have that
$$\underset{1h_0h}{}\underset{\begin{array}{c}\{1,\mathrm{},h\}\\ ||=k,h_0\end{array}}{}𝔖(\{h_0\})=(k+1)\underset{\begin{array}{c}\{1,\mathrm{},h\}\\ ||=k+1\end{array}}{}𝔖()\frac{h^{k+1}}{k!}.$$
Returning to the evaluation of $`\stackrel{~}{𝒮}_1`$, we find that
$$\stackrel{~}{𝒮}_1\left(\genfrac{}{}{0pt}{}{2\mathrm{}}{\mathrm{}}\right)\frac{N(\mathrm{log}N)(\mathrm{log}R)^{k+2\mathrm{}}h^k}{k!(k+2\mathrm{})!}\stackrel{~}{b}_1(k,\mathrm{},h)$$
where
$$\stackrel{~}{b}_1(k,\mathrm{},h)=2\frac{2\mathrm{}+1}{\mathrm{}+1}\frac{k}{k+2\mathrm{}+1}\frac{\mathrm{log}R}{\mathrm{log}N}+\frac{h}{\mathrm{log}N}1.$$
Unconditionally, we may take $`\theta =1/2`$, so $`\mathrm{log}R/\mathrm{log}N=1/4ϵ`$. We get two primes in some interval $`(n,n+h],N<n2N`$ provided $`\stackrel{~}{b}_1(k,\mathrm{},h)>0`$. This is equivalent to
$`{\displaystyle \frac{h}{\mathrm{log}N}}>`$ $`1{\displaystyle \frac{2k}{k+2\mathrm{}+1}}{\displaystyle \frac{2\mathrm{}+1}{\mathrm{}+1}}\left({\displaystyle \frac{1}{4}}ϵ\right)`$
$`=`$ $`{\displaystyle \frac{k+4\mathrm{}^2+6\mathrm{}+2+4ϵ(k+2k\mathrm{})}{2(1+\mathrm{})(1+2\mathrm{}+k)}}.`$
On letting $`\mathrm{}=[\sqrt{k}]`$ and taking $`k`$ sufficiently large, we see that this is valid with $`h/\mathrm{log}N`$ arbitrarily small. This proves Theorem 1.
For the proofs of Theorem 3 and Theorem 4, we note that if $`N<n2N`$ then
$$\varpi \varpi (n)\frac{(\mathrm{log}3N)^2}{2}.$$
Accordingly, we consider
(7.7) $`𝒮_2:={\displaystyle \underset{N<n2N}{}}`$ $`\{{\displaystyle \underset{h}{}}\varpi \varpi (n+h){\displaystyle \frac{(\mathrm{log}3N)^2}{2}}\}\times `$
$`\times \left({\displaystyle \underset{\mathrm{}=0}{\overset{L}{}}}b_{\mathrm{}}(\mathrm{log}R)^{\mathrm{}}\mathrm{\Lambda }_R(n;,\mathrm{})\right)^2.`$
The term $`n`$ contributes a negative amount unless there are two values $`h_i,h_j`$ such that $`n+h_i,n+h_j`$ are products of two primes. The values of $`n`$ for which any $`n+h`$ is a square of a prime contribute $`N^{1/2}(\mathrm{log}N)^{2k+2}`$, and this contribution may be absorbed into the error terms of our estimates.
We assume Hypotheses $`BV(\theta )`$ and $`BV_2(\theta )`$, and we argue along the same lines as in the proof of Theorem 2. When $`R=N^{(\theta ϵ)/2}`$, we obtain
$$𝒮_2=\underset{0\mathrm{}_1,\mathrm{}_2L}{}b_\mathrm{}_1b_\mathrm{}_2(\mathrm{log}R)^{\mathrm{}_1\mathrm{}_2}_2(\mathrm{}_1,\mathrm{}_2),$$
where
$`_2`$ $`𝔖()N(\mathrm{log}N)^2(\mathrm{log}R)^{k+\mathrm{}_1+\mathrm{}_2}(m_2(k,\mathrm{}_1,\mathrm{}_2,\theta )ϵ^{}),`$
$`m_2(k,\mathrm{}_1,\mathrm{}_2,\theta )`$ $`=m_{21}+m_{22}m_{23},`$
$`m_{21}`$ $`=\left({\displaystyle \genfrac{}{}{0pt}{}{\mathrm{}_1+\mathrm{}_2+2}{\mathrm{}_1+1}}\right){\displaystyle \frac{k}{(k+\mathrm{}_1+\mathrm{}_2+1)!}}{\displaystyle \frac{\theta }{2}},`$
$`m_{22}`$ $`=2\left\{\left({\displaystyle \genfrac{}{}{0pt}{}{\mathrm{}_1+\mathrm{}_2+2}{\mathrm{}_1+1}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{\mathrm{}_1+\mathrm{}_2+3}{\mathrm{}_1+1}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{\mathrm{}_1+\mathrm{}_2+3}{\mathrm{}_2+1}}\right)\right\}{\displaystyle \frac{k}{(k+\mathrm{}_1+\mathrm{}_2+2)!}}{\displaystyle \frac{\theta ^2}{4}},`$
$`m_{23}`$ $`={\displaystyle \frac{1}{2}}\left({\displaystyle \genfrac{}{}{0pt}{}{\mathrm{}_1+\mathrm{}_2}{\mathrm{}_1}}\right){\displaystyle \frac{1}{(k+\mathrm{}_1+\mathrm{}_2)!}},`$
and $`ϵ^{}=ϵ^{}(k,\mathrm{}_1,\mathrm{}_2,ϵ)0`$ as $`ϵ0`$.
Let $`𝐛`$ be as defined before. Then (suppressing the $`ϵ^{}`$ term)
$`𝒮_2^{}(N,,\theta ,𝐛)`$ $`:={\displaystyle \frac{𝒮_2}{N𝔖()(\mathrm{log}R)^k(\mathrm{log}N)^2}}{\displaystyle \underset{0\mathrm{}_1,\mathrm{}_2L}{}}b_\mathrm{}_1b_\mathrm{}_2m_2(k,\mathrm{}_1,\mathrm{}_2,\theta )`$
$`=𝐛^T𝐌_2𝐛,`$
where $`𝐌_2=𝐌_2(k,\theta )`$ is the matrix
$$𝐌_2=[m_2(k,i,j,\theta )]_{0i,jL}.$$
We first prove Theorem 4. As in the proof of Theorem 2, we wish to show that there is some $`𝐛`$ such that $`𝒮_2^{}>0`$ for a given $`\theta `$ and minimal $`k`$. Taking $`k=3`$ and $`L=1`$, we find that
(7.8)
$$𝐌_2=\frac{1}{480}\left[\begin{array}{cc}24\theta ^2+60\theta 40& 7\theta ^2+18\theta 10\\ 7\theta ^2+18\theta 10& 2\theta ^2+6\theta 4\end{array}\right].$$
If we take $`b(0)=1,b(1)=4`$, then we find that
$$𝐛^T𝐌_2𝐛=\frac{7\theta ^2}{30}+\frac{5\theta }{8}\frac{23}{60}.$$
This is positive whenever
$$\frac{75\sqrt{473}}{56}<\theta 1.$$
Finally, we note that $`=\{5,7,11\}`$ is an admissible $`3`$-tuple, so this completes the proof of Theorem 4.
We can also prove Theorem 4 with a slightly wider range of allowable $`\theta `$ by taking the determinant of the matrix in (7.8). A numerical calculation shows that this determinant has a zero at $`\theta =0.943635\mathrm{}`$.
For the proof of Theorem 3, we take $`k=8,L=2,\theta =1/2ϵ`$, and we find that
$$𝐌_2=\frac{1}{14!}\left[\begin{array}{ccc}216216& 8736& 3458\\ 8736& 364& 14\\ 3458& 14& 36\end{array}\right],$$
With
$$b(0)=1,b(1)=16,b(2)=16,$$
we find that
$$14!𝐛^T\mathrm{𝐌𝐛}=78760>0.$$
Now $`=\{11,13,17,19,23,29,31,37\}`$ is an admissible $`8`$-tuple, so this completes the proof of Theorem 3.
We make one final comment regarding the proofs that make use of bilinear forms in $`𝐛`$. By taking
$$\underset{\mathrm{}=0}{\overset{L}{}}b_{\mathrm{}}(\mathrm{log}R)^{\mathrm{}}\mathrm{\Lambda }_R(n;,\mathrm{})$$
in the definitions of $`𝒮_1`$ and $`𝒮_2`$, we are in essence using
$$y_r=𝔖()\underset{\mathrm{}=0}{\overset{L}{}}\frac{b_{\mathrm{}}}{\mathrm{}!}\left(\frac{\mathrm{log}R/r}{\mathrm{log}R}\right)^{\mathrm{}}.$$
In other words, we have essentially replaced $`(\mathrm{log}R/r)^{\mathrm{}}`$ in (1.20) by a polynomial in $`\mathrm{log}R/r`$.
Acknowledgements: We thank Tsz-Ho Chan and Yoichi Motohashi for their comments on this paper. Part of the work for this paper was done at the American Institute of Mathematics, where Graham was visiting in Fall 2004. He thanks them for their hospitality and excellent working environment. |
warning/0506/hep-th0506004.html | ar5iv | text | # Contents
## Chapter 1 Introduction
Gauge theories such as the Standard Model of particle physics have been investigated mainly in the framework of perturbation theory. Due to the tremendous complexity of (especially nonabelian) gauge theories, this is the only feasible approach for problems concerning single fundamental processes: It is not possible to solve even the simplest QED processes exactly. Concerning the macroscopic properties (e. g. thermodynamics) of a system, chances for a nonperturbative treatment in form of an effective theory are much better. Superconductivity in metals, for example, was described phenomenologically by the Landau-Ginzburg theory . Subsequently, a quantum theory of superconductivity was developed on the microscopic level .
In the Standard Model, a number of striking results has been obtained using perturbation theory, among them the anomalous magnetic moment of the electron, asymptotic freedom in QCD and many others. In spite of these great successes, perturbation theory has some intrinsic problems:
Perturbation theory is an expansion in powers of the gauge coupling. This is by definition only applicable to cases of small coupling. It is completely impossible to address strongly coupled problems in perturbation theory. Moreover, a perturbation expansion is an expansion about the trivial vacuum of the theory. But there exist objects in Yang-Mills theory, e. g. instantons, that are topologically distinct from the trivial vacuum. Such objects can by no means be included in an expansion about the trivial vacuum. This can also be seen from the partition function: The instanton enters the partition function with the measure $`e^S=e^{8\pi ^2/g^2}`$, which has an essential singularity at zero coupling such that the Taylor series of this function about $`g=0`$ vanishes identically. Therefore, excitations of nontrivial topology are completely ignored in perturbation theory.
For SU(N) Yang-Mills theory at finite temperature, further problems are known. It was shown that a perturbative calculation of the thermodynamical pressure can not be driven beyond order $`g^5`$ . This is essentially due to the existence of weakly screened, soft magnetic modes which cause infrared instabilities.
There is also a number of experimental and observational results in particle physics and cosmology the explanation of which is still open or being disputed. Examples are the nondetection of the Higgs particle at LEP, the existence of dark matter and dark energy in the universe, and a quadrupole signal in one of the power spectra of the cosmic microwave background. A nonperturbative description of gauge theories may be helpful or even necessary in understanding these aspects.
At present a large number of theories and approaches which try to include nonperturbative aspects in field theories is being considered. An analytical and nonperturbative approach to SU(2) and SU(3) Yang-Mills thermodynamics is presented in . In the spirit of the Ginzburg-Landau theory of superconductivity, a macroscopic field is used to account for microscopic processes in an effective theory. More precisely, a composite adjoint Higgs field $`\varphi `$, which describes the BPS saturated and topologically nontrivial part of the ground state, is introduced. The Higgs field $`\varphi `$ is generated at an asymptotically high temperature by noninteracting calorons of topological charge one and trivial holonomy. The field $`\varphi `$ is quantum mechanically and thermodynamically stabilized and can thus be used as a background for the topologically trivial sector of the theory. Interactions between trivial-holonomy calorons are included via a macroscopic pure-gauge configuration $`a_\mu ^{bg}`$ which is a solution to the equation of motion for the topologically trivial sector in the presence of the background $`\varphi `$.
As the modulus of the Higgs field decreases with temperature as $`|\varphi |\sqrt{\mathrm{\Lambda }^3/T}`$, where $`\mathrm{\Lambda }`$ is the Yang-Mills scale, the effects of topological defects die off at large temperature in a power-like fashion. Asymptotic freedom and infrared-ultraviolet decoupling of the fundamental theory, which are results obtained in perturbation theory at $`T=0`$, are preserved.
Thermodynamical quantities are in this framework calculated as loop expansions about the nontrivial ground state consisting of the Higgs field $`\varphi `$ and the pure-gauge configuration $`a_\mu ^{bg}`$. On tree level the excitations in the high temperature (or electric) phase are either massive thermal quasiparticles or massless ’photons’. The interactions between these quasiparticles appear to be very weak. The effective theory is both infrared- and ultraviolet-finite. The former is due to the existence of caloron-induced gauge boson masses (IR cutoff), the latter to constraints on loop momenta arising from the existence of the compositeness scale $`|\varphi |`$ (UV cutoff).
The purpose of this thesis is to compute the dynamical generation of the macroscopic, composite field $`\varphi `$ and to calculate the pressure of the Yang-Mills gas on two-loop level for the SU(2) case. The thesis is organized as follows: Chapter 2 reviews basic properties of Lie groups and pure Yang-Mills theory. The physics of some topological objects in gauge theories, namely the Abrikosov-Nielsen-Olesen vortex, the ’t Hooft-Polyakov monopole, and instantons is sketched. The focus is on the latter. Chapter 3 first presents a brief outline of the physics of the electric phase according to the approach in . The definition of the phase of the composite Higgs field in terms of a spatial and scale parameter average over an adjointly transforming two-point function is given and discussed. The average has to be evaluated on trivial holonomy caloron and anticaloron configurations. We discuss the uniqueness of the given definition and perform the evaluation. We show how under the assumption of an externally given scale the modulus of the field can be determined. The potential is deduced from the BPS equation. Chapter 4 contains the calculation of the two-loop contributions to the thermodynamical pressure of the SU(2) Yang Mills gas. We determine the contributing diagrams and state the Feynman rules. The computation of the diagrams is performed in the real time formalism of finite temperature field theory. The resulting integrals have to be evaluated numerically. We compare the two-loop contribution to the one-loop contribution and give an interpretation. Chapter 5 gives a short summary and an outlook on further research. The appendices contain technical details concerning the calculations in Chapters 3 and 4.
## Chapter 2 Basics of SU(N) Yang-Mills Theory and Solitonic Configurations
### 2.1 Lie groups and Lie algebras
For a continuous group $`G`$, the group elements may be parameterized by a set of continuously varying real parameters, $`\alpha ^a`$ ($`a=1,2,\mathrm{},K`$). The total number of parameters $`K`$ is called the order of the group. There are groups with compact parameter space such as SU(N), where $`K=\mathrm{N}^21`$, and SO(N), where $`K=\frac{\mathrm{N}(\mathrm{N}1)}{2}`$. There are also groups with non-compact parameter space, like the Poincaré group in four dimensions, where $`K=10`$.
Lie group. A Lie group is a continuous group $`G`$ where the set of parameters represents a differentiable manifold. The latter is referred to as group manifold. The multiplication map
$$G\times GG:(g_1,g_2)g_1g_2$$
(2.1)
and the inverse map
$$GG:gg^1$$
(2.2)
are differentiable. In terms of the parameters $`\alpha `$ and $`\alpha ^{}`$ this means
$$g(\alpha )g(\alpha ^{})=g(\alpha ^{\prime \prime }),$$
(2.3)
where $`\alpha ^{\prime \prime }`$ is an analytic function of $`\alpha `$ and $`\alpha ^{}`$ and similarly for the inverse map. As the dependence of the group elements on the parameters $`\alpha ^a`$ is analytic, any infinitesimal element $`gG`$ can be power expanded about the unit element of the group:
$$g(\alpha )=1+i\alpha ^aT^a+O(\alpha ^2).$$
(2.4)
The objects
$$T^a=\left(\frac{g(\alpha )}{\alpha ^a}\right)_{\alpha =0}$$
(2.5)
are called infinitesimal generators of the group $`G`$. For convenience, the unit element $`e`$ has parameters $`\alpha ^a=0`$, $`g(0)=e`$.
Lie algebra. A vector space $`V`$ together with a bilinear operation, the Lie-Bracket, given as
$$V\times VV:(X,Y)[X,Y]$$
(2.6)
and satisfying
$$[X,X]=0$$
(2.7)
and the Jacobi identity
$$[X,[Y,Z]]+[Y,[Z,X]]+[Z,[X,Y]]=0$$
(2.8)
is called a Lie algebra. The property Eq. (2.7) implies the antisymmetry of the Lie bracket,
$$[X,Y]=[Y,X].$$
(2.9)
The generators of a Lie group $`G`$ always form a Lie algebra $`𝔤`$. The dimension of the vector space $`V`$ is equal to the number of generators (that is the order) of $`G`$. If $`G`$ is a matrix group, then the Lie bracket is the usual matrix commutator and the Jacobi identity Eq. (2.8) is trivially fulfilled.
Since the generators $`T^a`$ of the group provide a basis of the Lie algebra, the Lie bracket of two generators must again be a linear combination of generators,
$$[T^a,T^b]=if^{abc}T^c.$$
(2.10)
The symbols $`f^{abc}`$ are called structure constants. They can be chosen to be completely antisymmetric and real. Using the structure constants, the Jacobi identity for the generators
$$[T^a,[T^b,T^c]]+[T^b,[T^c,T^a]]+[T^c,[T^a,T^b]]=0$$
(2.11)
can be phrased as
$$f^{ade}f^{bcd}+f^{bde}f^{cad}+f^{cde}f^{abd}=0.$$
(2.12)
Many properties of a Lie group can be derived from the Lie algebra; e. g. the commutation relations Eq. (2.10) of a Lie algebra (which themselves often are called the Lie algebra) completely determine the multiplication law of the associated Lie group in the vicinity of the unit element. The connection between Lie algebra and Lie group is established through the exponential map,
$$t𝔤\mathrm{exp}tG,$$
(2.13)
which, in case of a matrix group, is the usual exponential map defined by the power series
$$\mathrm{exp}A=\underset{n=0}{\overset{\mathrm{}}{}}\frac{A^n}{n!}.$$
(2.14)
There are Lie groups that have the same Lie algebra but group manifolds of different global structure and topology.
A Lie group $`G`$ that contains no invariant sub-group (ideal) other than $`\{e\}`$ and $`G`$ itself is called simple. If $`G`$ does not contain any Abelian invariant sub-group other than $`\{e\}`$, it is called semi-simple. The corresponding definitions apply to the Lie algebras.
For a Lie algebra, it is possible that some of the generators $`T^a`$ commute (i. e. their Lie-bracket is zero). These generators form the so-called Cartan sub-algebra, their number is known as the rank of the Lie algebra. For example SU(N) has rank $`\mathrm{N}1`$; SO(2) and SO(3) both have rank 1.
Representations. A linear representation $`R`$ of a group $`G`$ is a map which maps every group element $`gG`$ onto a linear transformation $`R(g)`$ of a vector space $`V`$. This map has to respect the group multiplication, i. e.
$$R(g_1)R(g_2)=R(g_1g_2)$$
(2.15)
and to map the unit element $`e`$ onto the unit matrix,
$$R(e)=\mathrm{𝟙}.$$
(2.16)
In other words, this map is a homomorphism. When a basis for $`V`$ is chosen, $`R(g)`$ are matrices. The dimension of the representation is by definition the dimension of the vector space $`V`$. In quantum mechanics, one is usually interested in finite dimensional unitary representations because unitarity is closely connected to probability density conservation.
A representation is called reducible if it is possible to find a basis in which all representation matrices have the form
$$R(g)=\left(\begin{array}{cc}R^1(g)& A(g)\\ 0& R^2(g)\end{array}\right),$$
(2.17)
otherwise it is called irreducible. If additionally $`A(g)=0`$ for all $`gG`$, the representation is called fully reducible. For a semi-simple group, all reducible representations are fully reducible (Weyl’s theorem).
Two representations $`R`$, $`R^{}`$ of a group are said to be equivalent, if they only differ by a similarity transformation
$$R^{}(g)=S^1R(g)SgG$$
(2.18)
with a nonsingular matrix $`S`$. It can be shown that every representation of a compact group is equivalent to a unitary representation.
To form a representation of the Lie algebra, the representation matrices have to respect the commutation relations Eq. (2.10) of the Lie algebra. The elements of a matrix group can be viewed as linear transformations of $`^n`$. Thus the group elements themselves form a linear representation of the group, the fundamental representation. The $`r`$ generators of the Lie group $`G`$ carry a representation of dimension $`r`$, the adjoint representation. The generator $`T^a`$ is mapped onto the mapping (denoted by the same symbol $`T^a`$)
$$T^a:𝔤𝔤;T^b[T^a,T^b].$$
(2.19)
The representation matrices of the adjoint representation are given by the structure constants,
$$R_{ac}^{\text{adj}}(T^b)=if^{abc}.$$
(2.20)
As a direct consequence of the Jacobi identity Eq. (2.8), the adjoint representation matrices fulfill the Lie algebra commutation relations Eq. (2.10), as requested. The adjoint representation is always a real representation because the structure constants are real and antisymmetric. The adjoint representation of a simple Lie group is always irreducible.
The group SU(N). SU(N) is the group of unitary $`\mathrm{N}\times \mathrm{N}`$ complex matrices with unit determinant,
$$U^{}U=\mathrm{𝟙},detU=1.$$
(2.21)
It is a compact simple Lie group. Its Lie algebra has rank $`\mathrm{N}1`$. The generators of SU(N) are the Hermitian and traceless $`\mathrm{N}\times \mathrm{N}`$ matrices
$$H^{}=H,\mathrm{tr}H=0.$$
(2.22)
The generators of SU(2) and SU(3) are usually taken to be the Pauli and Gell-Mann matrices, respectively. The adjoint representation of SU(N) has dimension $`\mathrm{N}^21`$.
### 2.2 SU(N) Yang-Mills theory
We are considering a four-dimensional Minkowskian spacetime. An SU(N) Yang-Mills theory is governed by a Lagrangian which is invariant under any local SU(N) transformation. Every field in the theory has to transform under a unitary and finite dimensional representation of SU(N). A Yang-Mills theory containing only gauge fields but no matter fields is often called pure.
The ordinary derivative $`_\mu `$ cannot be used to construct gauge invariant quantities. So, in order to be able to include derivative terms (i. e. kinetic terms for the fields) in the Lagrangian, the covariant derivative $`D_\mu `$ and a gauge field $`A_\mu `$ have to be introduced<sup>1</sup><sup>1</sup>1 One can conventionally absorb the gauge coupling $`e`$ in the definition of the gauge field and write $`D_\mu =_\mu iA_\mu `$ and $`F_{\mu \nu }=_\mu A_\nu _\nu A_\mu i[A_\mu ,A_\nu ]`$ etc. This notation is convenient for considering nonperturbative aspects. We will use it for working with instantons in Secs. 2.5, 2.6 and in Chapter 3.,
$$D_\mu =_\mu ieA_\mu .$$
(2.23)
The gauge field $`A_\mu `$ can be expanded in terms of the Hermitian generators of SU(N),
$$A_\mu =A_\mu ^a\frac{\lambda ^a}{2},$$
(2.24)
where the SU(N) generators $`\lambda ^a`$ are normalized such that $`\mathrm{tr}\lambda ^a\lambda ^b=2\delta _{ab}`$. When applying the covariant derivative to a matter field, the gauge field $`A_\mu `$ is understood to act in the representation of the matter field. That is for a fundamental field $`\phi `$
$$D_\mu \phi =_\mu \phi ieA_\mu ^a\frac{\lambda ^a}{2}\phi $$
(2.25)
and for an adjoint field $`\varphi `$
$$(D_\mu \varphi )_a=_\mu \varphi _a+ef^{abc}A_\mu ^b\varphi _c.$$
(2.26)
The covariant derivative $`D_\mu `$ is constructed such that the covariant derivative of a field has exactly the same transformation law as the field itself. To satisfy this request, the gauge field has to transform under an SU(N) gauge transformation $`U(x)`$ according to
$$A_\mu (x)U(x)A_\mu (x)U^{}(x)+\frac{i}{e}U(x)_\mu U^{}(x).$$
(2.27)
Viewing the gauge group (being a Lie group) as an analytic manifold, the gauge field $`A_\mu `$ is a connection on this manifold. The curvature of the manifold (in differential geometry) corresponds to the field strength tensor (in field theory), namely
$$F_{\mu \nu }=\frac{i}{e}[D_\mu ,D_\nu ]=_\mu A_\nu _\nu A_\mu ie[A_\mu ,A_\nu ].$$
(2.28)
In contrast to Abelian theories, the field strength in nonabelian theories is not a gauge invariant quantity but transforms under the adjoint representation as
$$F_{\mu \nu }(x)U(x)F_{\mu \nu }(x)U^{}(x).$$
(2.29)
The field strength can also be written in matrix notation,
$$F_{\mu \nu }=F_{\mu \nu }^a\frac{\lambda ^a}{2},$$
(2.30)
with the components
$$F_{\mu \nu }^a=_\mu A_\nu ^a_\nu A_\mu ^a+ef^{abc}A_\mu ^bA_\nu ^c,$$
(2.31)
where $`f^{abc}`$ are the structure constants of the gauge group. The kinetic term for the gauge field $`A_\mu `$ in the Lagrangian is
$$=\frac{1}{2}\mathrm{tr}F_{\mu \nu }F^{\mu \nu }=\frac{1}{4}F_{\mu \nu }^aF^{\mu \nu a}.$$
(2.32)
If one demands Lorentz invariance, gauge symmetry, renormalizability and CP-invariance, no further terms are allowed in pure Yang-Mills theory. Relaxing the demand for CP-invariance, an additional term proportional to $`F_{\mu \nu }^a\stackrel{~}{F}^{\mu \nu a}`$ may be added to the Lagrangian in Eq. (2.32). Here $`\stackrel{~}{F}_{\mu \nu }=\frac{1}{2}\epsilon _{\mu \nu \alpha \beta }F^{\alpha \beta }`$ denotes the dual field strength.
In particular a mass term for the gauge field is forbidden because it has the gauge variant form $`m^2\mathrm{tr}A_\mu A_\mu `$. Nevertheless, gauge bosons can acquire mass by dynamical symmetry breaking which is manifested by the Higgs mechanism.
The commutator term in Eq. (2.28) vanishes for Abelian gauge groups, so the Lagrangian in this case is quadratic in the gauge field, the equations of motion are linear in the gauge field, and hence there is no self interaction. If, in contrast, the gauge group is nonabelian (as is SU(N)), this is no longer true. The Lagrangian contains terms cubic and quartic in the gauge field (besides quadratic terms) and thus allows for three- and four-gauge boson vertices. Because of this a pure SU(N) Yang-Mills theory is interacting.
The equations of motion for the field $`A_\mu `$ are derived via the minimal action principle. For SU(N) pure Yang-Mills theory, they are
$$D_\mu F^{\mu \nu }=0$$
(2.33)
or in components
$$^\mu F_{\mu \nu }^a+ef^{abc}A^{b\mu }F_{\mu \nu }^c=0.$$
(2.34)
The right hand side of the equation of motion is zero because of the absence of external sources, i. e. charged matter. Nevertheless, this is an interacting theory since the gauge field couples to itself due to the nonlinearity in the field tensor.
On the classical level the Lagrangian Eq. (2.32) does not contain any dimensionful parameter: The gauge coupling $`e`$ is dimensionless, and gauge boson masses are forbidden. This is the reason for the invariance of the action under a rescaling of fields and spacetime arguments. On the quantum level, however, a mass scale comes into existence due to the mechanism of dimensional transmutation .
### 2.3 The Abrikosov-Nielsen-Olesen vortex line
Consider a theory with a U(1) gauge field $`A_\mu `$ and a scalar Higgs field $`\varphi `$ defined by the Lagrangian
$$=\frac{1}{4}F_{\mu \nu }F^{\mu \nu }+(D_\mu \varphi )(D^\mu \varphi )m^2\varphi ^{}\varphi \lambda (\varphi ^{}\varphi )^2.$$
(2.35)
If $`m^2<0`$, the gauge symmetry is spontaneously broken, and the Higgs field acquires a vacuum expectation value
$$|\varphi |_{\text{vac}}=\sqrt{\frac{m^2}{2\lambda }}.$$
(2.36)
The Higgs Lagrangian Eq. (2.35) exhibits static string-like solitonic solutions, so-called vortices. They were discovered by Nielsen and Olesen in 1973 . Consider a field configuration with cylindrical symmetry and the asymptotic behavior
$$\varphi =|\varphi |_{\text{vac}}e^{in\theta }\text{(}r\mathrm{}\text{)}$$
(2.37)
for the Higgs field, and
$$A_\mu =\frac{1}{e}_\mu (n\theta )\text{(}r\mathrm{}\text{)}$$
(2.38)
for the gauge field. In cylindrical coordinates $`z`$, $`r|𝒙|`$ and $`\theta `$, the latter reads
$$\begin{array}{cc}\hfill A_z& =0\hfill \\ \hfill A_r& =0\hfill \\ \hfill A_\theta & =\frac{n}{er}\text{(}r\mathrm{}\text{)}.\hfill \end{array}$$
(2.39)
Because of the demand for single valued fields, $`n`$ has to be an integer. In the asymptotic, $`A_\mu `$ is a pure gauge and thus the field strength vanishes,
$$F_{\mu \nu }=0\text{(}r\mathrm{}\text{)}.$$
(2.40)
The Higgs field $`\varphi `$ is covariantly constant,
$$D_r\varphi =0\text{and}D_\theta \varphi =0\text{(}r\mathrm{}\text{)},$$
(2.41)
and the potential $`V=m^2\varphi ^{}\varphi +\lambda (\varphi ^{}\varphi )^2`$ evaluates to zero. So the energy density at $`r\mathrm{}`$ is $`==0`$, and static solutions with finite energy and the above boundary conditions can exist in principal. The equations of motion obtained from the Langrangian Eq. (2.35) are
$$\begin{array}{cc}& D^\mu (D_\mu \varphi )=m^22\lambda \varphi |\varphi |^2\hfill \\ & ie(\varphi _\mu \varphi ^{}\varphi ^{}_\mu \varphi )+2e^2A_\mu |\varphi |^2=^\nu F_{\mu \nu }.\hfill \end{array}$$
(2.42)
One can check that the equations of motion allow for configurations with the asymptotic behavior Eqs. (2.37) and (2.38). To find a solution to the equations of motion which satisfies the above boundary conditions, one makes the ansatz
$$\begin{array}{cc}\hfill A_z(r)& =0\hfill \\ \hfill A_r(r)& =0\hfill \\ \hfill A_\theta (r)& A(r)\hfill \end{array}$$
(2.43)
for the gauge field, and
$$\varphi =\chi (r)e^{in\theta }$$
(2.44)
with
$$\begin{array}{c}\hfill \chi (r)\underset{r0}{\overset{}{}}0\text{and}\chi (r)\underset{r\mathrm{}}{\overset{}{}}|\varphi |_{\text{vac}}\end{array}$$
(2.45)
for the Higgs field. The magnetic field $`𝑩`$ will have only a $`z`$-component,
$$B_z=\frac{1}{r}\frac{d}{dr}[rA(r)].$$
(2.46)
Inserting the above ansatz into the equations of motion yields a system of differential equations for the functions $`A(r)`$ and $`\chi (r)`$,
$$\begin{array}{cc}& \frac{1}{r}\frac{d}{dr}\left(r\frac{d\chi }{dr}\right)\left[\left(\frac{n}{r}eA\right)^2+m^2+2\lambda \chi ^2\right]\chi =0\hfill \\ & \frac{d}{dr}\left(\frac{1}{r}\frac{d}{dr}(rA)\right)2e\left(\frac{n}{e}+eA\right)\chi ^2=0.\hfill \end{array}$$
(2.47)
No exact solutions to these equations are known. The asymptotic behavior of the gauge field and the magnetic field has been deduced by Nielsen and Olesen as
$$A=\frac{n}{er}\frac{c}{e}K_1\left(|e||\varphi |_{\text{vac}}r\right)\underset{r\mathrm{}}{\overset{}{}}\frac{n}{er}\frac{c}{e}\left(\frac{\pi }{2|e||\varphi |_{\text{vac}}r}\right)^{1/2}e^{|e||\varphi |_{\text{vac}}r}+\mathrm{}$$
(2.48)
and
$$B_z=c\chi K_0\left(|e||\varphi |_{\text{vac}}r\right)\underset{r\mathrm{}}{\overset{}{}}\frac{c}{e}\left(\frac{\pi |\varphi |_{\text{vac}}}{2|e|r}\right)^{1/2}e^{|e||\varphi |_{\text{vac}}r}+\mathrm{},$$
(2.49)
where $`K_0`$ and $`K_1`$ denote modified Bessel functions, and $`c`$ is a constant of integration.
The line integral over $`A_\mu `$ around a circle $`S_1`$ at infinity yields the magnetic flux through the surface enclosed,
$$\mathrm{\Phi }=𝑩𝑑𝝈=A_\mu 𝑑x^\mu =A_\theta r𝑑\theta =\frac{2\pi }{e}n,$$
(2.50)
The magnetic flux is quantized: it appears only in multiples of the flux quantum $`\frac{2\pi }{e}`$.
The vortex line owes its existence to the topological structure of the gauge manifold. The boundary condition Eq. (2.37) defines a mapping of the boundary $`S_1`$ in physical space onto the group manifold of U(1), which again is $`S_1`$. There are infinitely many classes of such mappings which can not be continuously deformed into one another,
$$\pi _1(U(1))=\pi _1(S_1)=.$$
(2.51)
A gauge theory, where the gauge group $`G`$ has $`\pi _1(G)=0`$, does not exhibit vortex lines. This is the case for SU(2), for example.
The Lagrangian Eq. (2.35) is the relativistic generalization of the Landau-Ginzburg free energy of a type II superconductor. The electromagnetic field $`A_\mu `$ interacts with the bosonic field $`\varphi `$, which describes the Cooper pairs. In the superconducting phase the photon becomes massive, and if an external magnetic field can enter the superconducting material (i. e. a type II superconductor), it does so only in the form of (quantized) Abrikosov flux tubes.
### 2.4 The ’t Hooft-Polyakov monopole
Magnetic monopoles in gauge theories were first considered in 1974 by ’t Hooft and Polyakov in the context of an SU(2) Yang-Mills theory with an isovector Higgs field $`\varphi ^a`$. The Langrangian of such a theory is
$$=\frac{1}{4}F_{\mu \nu }^aF^{\mu \nu a}+\frac{1}{2}(D_\mu \varphi ^a)(D^\mu \varphi ^a)\frac{m^2}{2}\varphi ^a\varphi ^a\lambda (\varphi ^a\varphi ^a)^2$$
(2.52)
with the parameter $`m^2`$ chosen negative such that the Higgs field has a non-zero vacuum expectation value $`F`$ with
$$F^2=\frac{m^2}{4\lambda }.$$
(2.53)
Now, consider static solutions with the asymptotic behavior
$$\begin{array}{cc}\hfill A_i^a& =\epsilon _{iab}\frac{x_b}{er^2}\text{(}r\mathrm{}\text{)}\hfill \\ \hfill A_0^a& =0\hfill \end{array}$$
(2.54)
for the SU(2) gauge field and
$$\varphi ^a=F\frac{x^a}{r}\text{(}r\mathrm{}\text{)}$$
(2.55)
for the Higgs field. Here $`r`$ is the norm of the spatial vector, $`r=|𝒙|`$. Note that in Eqs. (2.54) and (2.55) space and isospace indices are mixed. Fields of the form Eq. (2.55) are known as ”hedgehogs”.
By the presence of a nonvanishing vacuum expectation value of the Higgs triplet the SU(2) gauge symmetry is broken. The field strength corresponding to the unbroken U(1) subgroup is the ’t Hooft tensor
$$F_{\mu \nu }=\frac{1}{|\varphi |}\varphi ^aF_{\mu \nu }^a\frac{1}{e|\varphi |^3}\epsilon _{abc}\varphi ^a(D_\mu \varphi ^b)(D_\nu \varphi ^c).$$
(2.56)
Upon inserting
$$\begin{array}{cc}\hfill A_\mu ^{1,2}=0& A_\mu ^3=A_\mu \hfill \\ \hfill \varphi ^{1,2}=0& \varphi ^3=F,\hfill \end{array}$$
(2.57)
it reduces to the usual definition of the electro-magnetic field tensor. Inserting the nontrivial asymptotic conditions Eqs. (2.54) and (2.55) the ’t Hooft tensor evaluates to
$$F_{\mu \nu }=\frac{1}{er^3}\epsilon _{\mu \nu a}x^a\text{(}r\mathrm{}\text{)}.$$
(2.58)
This corresponds to the radial magnetic field of magnetic point charge,
$$B_a=\frac{x^a}{er^3}\text{(}r\mathrm{}\text{)}$$
(2.59)
with a total flux or magnetic charge
$$q_{mag}=\mathrm{\Phi }=_{S_2}𝑑\sigma _aB_a=\frac{4\pi }{e}.$$
(2.60)
In it is shown that there exist configurations with the requested asymptotic behavior that are smooth and hence have finite energy for all $`\lambda `$ and $`m^2<0`$. In they are explicitly given for the so-called BPS limit $`\lambda 0`$. In the BPS limit, the mass of the monopole is
$$M_\text{m}=\frac{4\pi }{e^2}M_\text{W},$$
(2.61)
where $`M_W=eF=\frac{em}{2\sqrt{\lambda }}`$ is the vector boson mass. For general $`\lambda `$, the mass of the monopole is larger but still of the same order.
The magnetic charge can also be expressed as
$$q_{mag}=\frac{1}{4\pi }d^3xK^0$$
(2.62)
with the current
$$\begin{array}{cc}\hfill K^\mu & =\frac{1}{2e}\epsilon ^{\mu \nu \rho \sigma }\epsilon _{abc}_\nu \widehat{\varphi }^a_\rho \widehat{\varphi }^b_\sigma \widehat{\varphi }^c\hfill \\ & =_\nu \stackrel{~}{F}_{\mu \nu },\hfill \end{array}$$
(2.63)
where $`\widehat{\varphi }^a=\frac{\varphi ^a}{|\varphi |}`$. This current is identically conserved, $`_\mu K^\mu =0`$; it is not a Noether current corresponding to some symmetry of the Lagrangian, but of a topological nature. Inserting Eq. (2.63) into Eq. (2.62) and applying Gauss’ theorem yields
$$q_m=\frac{1}{8\pi e}_{S_2}𝑑\sigma _i\epsilon _{ijk}\epsilon _{abc}\widehat{\varphi }^a_j\widehat{\varphi }^b_k\widehat{\varphi }^c,$$
(2.64)
where the integral has to be performed over an $`S_2`$ with infinite radius. From Eq. (2.64) it can be seen that the total magnetic flux is completely carried by the Higgs field. Moreover, Eq. (2.64) states that the magnetic charge does only depend on the asymptotic behavior of the fields. The existence of configurations with non-zero magnetic charge is due to the possibility to demand nontrivial boundary conditions for the fields, that means: There are maps from the surface of space ($`S_2`$) onto the manifold $`S_2`$ (corresponding to rotations in isospace) which can not be continuously deformed into the constant map. Mathematically speaking, the second homotopy group of $`S_2`$ is nontrivial,
$$\pi _2(S_2)=.$$
(2.65)
Conservation of magnetic monopole charge is due to this topological argument, namely the transition between nonhomotopic gauge configurations needs infinite energy.
### 2.5 Instantons at zero temperature
Introductory material on instantons and other topological objects in gauge theories is presented in . Reviews on instanton physics are .
SU(N) vacuum. The (Minkowskian) vacuum of an SU(N) Yang-Mills theory is infinitely degenerate. There is not only the trivial vacuum $`A_\mu (x)0`$, but every field configuration of the form
$$A_\mu (x)=iU(x)_\mu U^{}(x),$$
(2.66)
where $`U(x)`$ is an SU(N) matrix, differs from $`A_\mu =0`$ only by a gauge transformation and hence has zero field strength and energy density, as well. Configurations as in Eq. (2.66) are referred to as pure gauge. These pure gauge configurations fall into topologically distinct classes, and hence they cannot be smoothly connected. The configurations Eq. (2.66) are classified by a winding number, the Pontryagin index, defined as
$$n_W=\frac{1}{24\pi ^2}d^3x\epsilon ^{ijk}\mathrm{tr}(U^{}_iU)(U^{}_jU)(U^{}_kU).$$
(2.67)
The quantity $`n_W`$ is an invariant under smooth deformations. It can also be expressed in terms of the gauge field as
$$n_W=\frac{1}{16\pi ^2}d^3x\epsilon ^{ijk}\left(A_i^a_jA_k^a+\frac{1}{3}f^{abc}A_i^aA_j^bA_k^c\right).$$
(2.68)
It is not possible to find a solution to the equation of motion with finite energy that connects two vacua with different winding number. Configurations corresponding to a transition between different vacua of the theory via a tunneling process are called instantons. They have to be considered in the framework of Euclidean field theory.
Tunneling solutions. Tunneling configurations are solutions to the Euclidean equations of motion (i. e. they minimize the action of the theory, formulated in a Euclidean spacetime). This can be motivated by an example from the quantum mechanics of a point particle.
Imagine a point particle with kinetic energy $`T`$ moving in a time independent potential $`V`$. A semiclassical approximation (that is an expansion around the classical path) will not include any tunneling process because there are no classical tunneling paths. Now letting
$$EVVE$$
(2.69)
clearly interchanges allowed and forbidden region, so that a classical path exists, see Fig. 2.1. An expansion around this classical path yields the WKB approximation for the tunneling process. But $`EV=T`$, and $`T`$ is the square of a time derivative. Hence the sign change Eq. (2.69) is equivalent to the Euclidean continuation,
$$ti\tau \text{(}\tau \text{ real)}.$$
(2.70)
So tunneling processes are solutions to the Euclidean equations of motion. This is true in quantum field theory as well.
Finite action and boundary conditions. In order to get finite action configurations, one has to demand that the field strength $`F_{\mu \nu }`$ approaches zero at the boundary of spacetime faster than $`1/x^2`$. This can not only be fulfilled trivially (by approaching $`A_\mu =0`$), but also nontrivially by demanding that the gauge field $`A_\mu `$ is pure gauge on the boundary of Euclidean spacetime (which is $`S_3`$),
$$A_\mu =iU(x)_\mu U^{}(x)\text{(}\left|x\right|\mathrm{}\text{)}.$$
(2.71)
This condition defines a mapping from the surface $`S_3`$ to the group space of SU(N). The equivalence classes of homotopic mappings from $`S_3`$ into a manifold $`X`$, $`f:S_3X`$, form the so-called third homotopy group of the manifold, $`\pi _3(X)`$. So the finite action field configurations fall into topologically distinguished classes if and only if the third homotopy group of the gauge group manifold is nontrivial. The homotopy group $`\pi _3`$ of most gauge groups is known, see for example . The case of basic interest is SU(2), where
$$\pi _3(\text{SU(2)})=\pi _3(S_3)=.$$
(2.72)
The integer attached to a given mapping $`f:S_3S_3`$ indicates how often the gauge group manifold $`S_3`$ is ”wrapped” around the spacetime-$`S_3`$ and hence is referred to as winding number. It is also known as the Brouwer degree of the mapping $`f`$.
Pontryagin index and topological charge. For a given field configuration $`F_{\mu \nu }`$ the four dimensional Pontryagin index is defined as
$$Q=\frac{1}{32\pi ^2}d^4xF_{\mu \nu }^a\stackrel{~}{F}_{\mu \nu }^a.$$
(2.73)
The integrand of Eq. (2.73), the Pontryagin density, can also be written as a total divergence,
$$\frac{1}{32\pi ^2}F_{\mu \nu }^a\stackrel{~}{F}_{\mu \nu }^a=_\mu K_\mu ,$$
(2.74)
where $`K_\mu `$ is the Chern-Simons current,
$$K_\mu =\frac{1}{16\pi ^2}\epsilon _{\mu \alpha \beta \gamma }\left(A_\alpha ^a_\beta A_\gamma ^a+\frac{1}{3}f^{abc}A_\alpha ^aA_\beta ^bA_\gamma ^c\right).$$
(2.75)
In contrast to the Pontryagin density, the Chern-Simons current is a gauge variant quantity. If the integrand is nonsingular, the volume integral in Eq. (2.73) can be converted to a surface integral by means of Gauss’ theorem,
$$Q=d^4x_\mu K_\mu =𝑑\sigma _\mu K_\mu .$$
(2.76)
Although this surface integral has to be evaluated on a sphere $`S_3`$ with infinite radius, it is not necessarily zero: If the gauge field $`A_\mu `$ falls off less than $`1/|x|`$, then the Pontryagin index will be non-zero.
From Eq. (2.76) it is obvious that (for regular gauge fields) the Pontryagin index is solely determined by the asymptotic behavior of the gauge field or in other words by the boundary conditions. Thus the existence of fields with non-zero Pontryagin index is due to the possibility of demanding nontrivial boundary conditions, as already indicated in the context of Eq. (2.72). The Pontryagin index is a topological charge because its conservation does not follow from a continuous symmetry of the Lagrangian but is solely due to the stability of the boundary conditions under continuous perturbations.
Bogomolnyi bound and self-duality. For the construction of a configuration of minimal Euclidean action connecting two vacua with different winding (Chern-Simons) number, the Yang-Mills action is written in form of a Bogomolnyi decomposition,
$$\begin{array}{cc}\hfill S& =\frac{1}{4g^2}d^4xF_{\mu \nu }^aF_{\mu \nu }^a\hfill \\ & =\frac{1}{4g^2}d^4x\left(\pm F_{\mu \nu }^a\stackrel{~}{F}_{\mu \nu }^a+\frac{1}{2}(F_{\mu \nu }^a\stackrel{~}{F}_{\mu \nu }^a)^2\right)\hfill \end{array}$$
(2.77)
where
$$\stackrel{~}{F}_{\mu \nu }=\frac{1}{2}\epsilon _{\mu \nu \alpha \beta }F_{\alpha \beta }$$
(2.78)
is the dual field strength tensor. The second term in Eq. (2.77) is a square and hence always positive; the first term is the Pontryagin index $`Q`$ defined as in Eq. (2.73), a topological invariant. Hence any self-dual (or anti self-dual) field configuration, i. e.
$$F_{\mu \nu }=\pm \stackrel{~}{F}_{\mu \nu },$$
(2.79)
is a minimum of the Euclidean action (in a given topological sector). It has the action
$$S=\frac{8\pi ^2}{g^2}|Q|.$$
(2.80)
In the dual field strength $`\stackrel{~}{F}_{\mu \nu }`$ the roles of electric and magnetic fields are interchanged as opposed to $`F_{\mu \nu }`$. So, (anti) self-duality can (at least in temporal gauge) also be characterized by
$$𝑬^a=\pm 𝑩^a,$$
(2.81)
where $`𝑬^a`$ and $`𝑩^a`$ are the color electric and magnetic fields respectively. From the decomposition Eq. (2.77) it is clear that a self-dual gauge field is a (local) minimum of the Yang-Mills action, and indeed self-duality and the Bianchi identity imply that the equations of motion are satisfied,
$$D_\mu F_{\mu \nu }=\pm D_\mu \stackrel{~}{F}_{\mu \nu }=0.$$
(2.82)
In contrast to the equation of motion (2.82), which is a second-order differential equation, the self-duality equation (2.79) is first order. The Bogomolnyi decomposition Eq. (2.77) gives a lower bound for the action in a given topological sector. Configurations saturating this bound (i. e. self-dual configurations) are also referred to as BPS-saturated, . It can be shown that the energy-momentum tensor vanishes identically on BPS-saturated fields.
BPST instanton. Here we consider the gauge group SU(2). In the boundary condition
$$A_\mu iU_\mu U^{}$$
(2.83)
which is necessary to achieve a configuration of finite Euclidean action (see above), one chooses
$$U(x)=\frac{x_4+ix_i\lambda _i}{|x|}\text{(}|x|\mathrm{}\text{)},$$
(2.84)
where $`|x|=\sqrt{x_\mu x_\mu }`$, and $`\lambda _i`$ the are Pauli matrices. Therefore the gauge potential has to have the asymptotic form
$$A_\mu ^a2\eta _{\mu \nu }^a\frac{x_\nu }{x^2}\text{(}|x|\mathrm{}\text{)},$$
(2.85)
where the ’t Hooft symbols $`\eta _{\mu \nu }^a`$ and $`\overline{\eta }_{\mu \nu }^a`$ are defined as
$$\begin{array}{cc}\hfill \eta _{\mu \nu }^a& =\epsilon _{a\mu \nu }+\delta _{a\mu }\delta _{\nu 4}\delta _{a\nu }\delta _{\mu 4}\hfill \\ \hfill \overline{\eta }_{\mu \nu }^a& =\epsilon _{a\mu \nu }\delta _{a\mu }\delta _{\nu 4}+\delta _{a\nu }\delta _{\mu 4}.\hfill \end{array}$$
(2.86)
This leads to the ansatz
$$A_\mu ^a=2\eta _{\mu \nu }^a\frac{x_\nu f(x^2)}{x^2},$$
(2.87)
where the scalar function $`f`$ has to satisfy the boundary condition $`f1\text{ for }|x|\mathrm{}`$. Inserting this ansatz into the self-duality equation (2.79) yields the differential equation
$$f(1f)x^2f^{}=0,$$
(2.88)
which is solved by
$$f=\frac{x^2}{x^2+\rho ^2},$$
(2.89)
where $`\rho `$ is a constant of integration. This gauge field configuration is the Belavin-Polyakov-Schwartz-Tyupkin instanton
$$A_\mu ^a(x)=2\eta _{\mu \nu }^a\frac{x_\nu }{x^2+\rho ^2}.$$
(2.90)
For this instanton solution, the Chern-Simons current Eq. (2.75) has the asymptotic behavior
$$K_\mu \frac{1}{2\pi ^2}\frac{x_\mu }{|x|^4}\text{(}|x|\mathrm{}\text{)},$$
(2.91)
and inserting this into Eq. (2.76) yields $`Q=1`$ for the topological charge of the BPST instanton. Replacing $`\eta _{\mu \nu }^a`$ with $`\overline{\eta }_{\mu \nu }^a`$ yields a solution with $`Q=1`$. The corresponding field strength is
$$F_{\mu \nu }^a=4\eta _{\mu \nu }^a\frac{\rho ^2}{(x^2+\rho ^2)^2}.$$
(2.92)
Although the instanton gauge potential only falls of as $`1/x`$, the field strength $`F_{\mu \nu }^a`$ falls of as $`1/x^4`$. So instantons are well localized in (Euclidean) time and in space.
The solution Eq. (2.90) can be generalized by shifting the instanton center to an arbitrary point $`z_\mu `$,
$$A_\mu ^a(x)=2\eta _{\mu \nu }^a\frac{(xz)_\nu }{(xz)^2+\rho ^2},$$
(2.93)
and changing the color orientation of the instanton by a global gauge rotation (containing three parameters). Furthermore the BPST instanton contains an arbitrary parameter $`\rho `$, which gives the size of the instanton. These make up eight parameters in the generalized BPST instanton. Since the action of the configuration Eq. (2.93), namely
$$S=\frac{8\pi ^2}{g^2}\left|Q\right|=\frac{8\pi ^2}{g^2},$$
(2.94)
is independent of these eight parameters, they correspond to zero modes. The gauge in which the BPST instanton has the form (2.90) is referred to as regular gauge.
Singular gauge. The BPST instanton Eq. (2.90) has a singularity at infinity. This singularity can by means of the singular gauge transformation
$$U(x)=\frac{x_4+ix_i\lambda _i}{|x|}$$
(2.95)
be shifted from infinity to the origin, such that the instanton in so-called singular gauge has the gauge potential
$$A_\mu ^a(x)=2\overline{\eta }_{\mu \nu }^a\frac{x_\nu }{x^2}\frac{\rho ^2}{x^2+\rho ^2}=\overline{\eta }_{\mu \nu }^a_\nu \mathrm{ln}\left(1+\frac{\rho ^2}{(xz)^2}\right).$$
(2.96)
The singularity of the gauge field $`A_\mu `$ at $`x=0`$ is not physical; field strength, action density and topological charge density are smooth. The field strength of an instanton in singular gauge is
$$F_{\mu \nu }^a(x)=\frac{4\rho ^2}{(x^2+\rho ^2)^2}\left(\overline{\eta }_{\mu \nu }^a2\overline{\eta }_{\mu \kappa }^a\frac{x_\kappa x_\nu }{x^2}2\overline{\eta }_{\kappa \nu }^a\frac{x_\mu x_\kappa }{x^2}\right).$$
(2.97)
In regular gauge the integral over the topological charge density picked up only contributions at $`|x|\mathrm{}`$, cf. Eq. (2.76). In contrast to that, in singular gauge not only the action density, but also the topological charge density is localized at the center of the instanton. It is this property that allows for the construction of multi instantons in singular gauge.
Multi instanton solutions. The ’t Hooft ansatz for the gauge potential
$$A_\mu ^a(x)=\overline{\eta }_{\mu \nu }^a_\nu \mathrm{ln}\mathrm{\Pi }(x)$$
(2.98)
is a generalization of the BPST instanton in singular gauge, Eq. (2.96). The demand for self-duality translates into a condition for the scalar pre-potential $`\mathrm{\Pi }(x)`$, it has to fulfill the Laplace equation
$$\mathrm{\Pi }^1(x)_\mu _\mu \mathrm{\Pi }(x)=0.$$
(2.99)
Hence there is (only in singular gauge) some kind of superposition principle: Superposing the pre-potentials of single instantons with different centers and scale parameters again yields a self-dual solution to the equations of motion, a multi instanton. Therefore the pre-potential $`\mathrm{\Pi }(x)`$ has the general form
$$\mathrm{\Pi }(x)=1+\underset{n=1}{\overset{K}{}}\frac{\rho _n^2}{(xz_n)^2}.$$
(2.100)
The configuration Eq. (2.98) with the pre-potential Eq. (2.100) inserted has Pontryagin index $`Q=K`$. Replacing $`\eta _{\mu \nu }^a`$ in the ’t Hooft ansatz with $`\overline{\eta }_{\mu \nu }^a`$ yields a solution with $`Q=K`$. The ’t Hooft ansatz in particular allows for the construction of periodic instantons.
Collective coordinates. The collective coordinates of the SU(2) instanton with topological charge $`|Q|=1`$ are the instanton size $`\rho `$, the instanton position $`z`$ and three parameters determining the color orientation, making up 8 parameters in total. The action on an instanton is independent of all theses parameters; i. e. they are moduli of the instanton solution and correspond to zero modes.
A multi instanton of topological charge $`K`$ is the superposition of $`K`$ single instantons and accordingly has $`8K`$ parameters, . Note that the multi instanton solution as displayed in Eqs. (2.98) and (2.100) does not cover all of this moduli space: the $`K`$ instantons are chosen to have the same color orientation, which must not necessarily be the case. The general SU(2) multi instanton solution has been constructed by Atiyah et. al. .
Gauge groups other than SU(2). For the case of a simple gauge group $`G`$, basically no instantons other than in SU(2) arise, but the various possible embeddings of SU(2) in $`G`$ have to be taken into account. The general SU(N) multi-instanton with Pontryagin index $`Q`$ has $`4\mathrm{N}Q`$ parameters, .
Tunneling interpretation. Assuming that the (regular) gauge potential falls of rapidly at spatial infinity, one can write the topological charge $`Q`$ from Eq. (2.76) as
$$\begin{array}{cc}\hfill Q& =d^4x_\mu K_\mu \hfill \\ & =𝑑\tau \frac{d}{d\tau }d^3xK_0+𝑑\tau 𝑑\sigma _iK_i\hfill \\ & =_{\tau =\mathrm{}}d^3xK_0+_{\tau =\mathrm{}}d^3xK_0\hfill \\ & =n_W(\tau =\mathrm{})n_W(\tau =\mathrm{}),\hfill \end{array}$$
(2.101)
where the fact that the zero component of the Chern-Simons current Eq. (2.75) indeed is equal to the integrand of Eq. (2.68) was used. That means an instanton with Pontryagin index $`Q0`$ connects topologically different vacua with the according difference in winding number via tunneling.
### 2.6 Instantons at finite temperature
Finite temperature. The equilibrium thermodynamics of a quantum field theory at finite temperature usually is considered in the framework of Euclidean field theory. The finite temperature $`T`$ corresponds to imaginary time being compactified on a circle with circumference $`\beta =\frac{1}{T}`$. The physical fields have to be periodic in Euclidean time. For a gauge field the periodicity condition is less strict; only periodicity up to a gauge transformation is demanded,
$$A_\mu (\beta ,𝒙)=\mathrm{\Omega }(0,𝒙)A_\mu (0,𝒙)\mathrm{\Omega }^{}(0,𝒙)+\frac{i}{g}\mathrm{\Omega }(0,𝒙)_\mu \mathrm{\Omega }^{}(0,𝒙)$$
(2.102)
with $`\mathrm{\Omega }\text{SU(N)}`$. We now have to discern time and space and write $`x_\mu =(\tau ,𝒙)`$.
Periodic instanton. To find an instanton solution which is periodic in Euclidean time, one starts from the ’t Hooft ansatz Eq. (2.98) with the pre-potential Eq. (2.100) describing a multi instanton built out of $`K`$ instantons with scales $`\rho _n`$ and centered at $`z_n`$ respectively. All of them have the same color orientation. Taking the sum over infinitely many instanton pre-potentials centered at $`(n\beta ,0)`$, $`n`$ with the same scale parameter $`\rho `$ yields a periodic instanton (or caloron) centered at $`(0,0)`$,
$$\mathrm{\Pi }(\tau ,𝒙)=1+\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}\frac{\rho ^2}{(\tau n\beta ,𝒙)^2}.$$
(2.103)
The sum has been performed by Harrington and Shepard, see ,
$$\mathrm{\Pi }(\tau ,𝒙)=1+\frac{\pi \rho ^2}{\beta r}\frac{\mathrm{sinh}\frac{2\pi }{\beta }r}{\mathrm{cosh}\frac{2\pi }{\beta }r\mathrm{cos}\frac{2\pi }{\beta }\tau },$$
(2.104)
where $`r=\left|𝒙\right|`$. The caloron of Eq. (2.104) still has the action $`S=\frac{8\pi ^2}{g^2}`$; the action is, on the classical level, independent of temperature. By superposing (the pre-potentials of) single-calorons a caloron solution of higher topological charge can be constructed. As in the case of the multi-instanton, this is only possible in singular gauge. Calorons of higher topological charge also possess more moduli, some of which (namely the core sizes and the core distances) are dimensionful.
Polyakov loop. The Polyakov loop is a time-like Wilson loop
$$𝐏(𝒙)=𝒫\mathrm{exp}\left[ig_0^{1/T}𝑑\tau A_4(\tau ,𝒙)\right],$$
(2.105)
where $`𝒫`$ denotes the path ordering operation. It has to be evaluated in periodic gauge, $`A_\mu (\tau +\beta ,𝒙)=A_\mu (\tau ,𝒙)`$. The Polyakov loop is only defined at finite temperature. At spatial infinity, the Polyakov loop does not depend on $`𝒙/r`$ any longer. Its value $`𝐏(|𝒙|\mathrm{})`$ is a topological invariant.
Calorons are classified according to the eigenvalues of their Polyakov loop at spatial infinity. By definition, trivial holonomy means for SU(2) calorons that $`𝐏(|𝒙|\mathrm{})=\pm \mathrm{𝟙}`$. Calorons with non-unity eigenvalue of $`𝐏`$ at spatial infinity are said to have nontrivial holonomy, accordingly. The caloron solution from Eq. (2.104) has Polyakov loop $`𝐏=\mathrm{𝟙}`$. More general solutions have been constructed by Nahm , Lee, Lu , and Kraan, van Baal
Monopole constituents. In 1998 it has been shown independently by Lee and Lu and Kraan and van Baal that SU(N) calorons of nontrivial holonomy contain constituent BPS magnetic monopoles. These monopoles are subject to an attractive interaction in the case of small holonomy (i. e. close to trivial holonomy), or to a repulsive potential for large holonomy (i. e. far from trivial holonomy), see . The latter can lead to dissociation of the caloron into a monopole-antimonopole pair.
## Chapter 3 Composite Adjoint Higgs Field in SU(2) Yang-Mills Theory
In an analytical and nonperturbative approach to the thermodynamics of SU(2) and SU(3) Yang-Mills theory in four dimensions is developed. An essential ingredient is the existence of an adjoint Higgs field $`\varphi `$ composed of trivial holonomy calorons.
Sec. 3.1 gives a brief outline of this approach as presented in . We confine ourselves to the case of SU(2) and its so-called electric phase. We will give and evaluate a microscopic definition of the composite field $`\varphi `$ in Sec. 3.2, and briefly comment on the nonexistence of alternative possibilities for its definition.
### 3.1 An outline of the approach in
Decomposition of gauge fields. Each SU(2) gauge field $`A_\mu `$ appearing in the partition function of the fundamental theory is uniquely decomposed into a topologically nontrivial and BPS saturated part $`A_\mu ^{top}`$ represented by calorons, and a topologically trivial remainder $`a_\mu `$,
$$A_\mu =A_\mu ^{top}+a_\mu .$$
(3.1)
At a large temperature, calorons with topological charge one and trivial-holonomy are assumed to generate a macroscopic adjoint scalar field $`\varphi `$. Interactions between the trivial holonomy calorons via topologically trivial fluctuations are not included in the field $`\varphi `$ but are accounted for at a later stage by means of a pure-gauge background $`a_\mu ^{bg}`$. The field $`\varphi `$ has to have the following properties:
1. It describes (part of) the ground state of a thermal system, so its (gauge invariant) modulus $`|\varphi |`$ must be independent of space and time.
2. The action of a classical caloron is independent of temperature. So no explicit $`T`$ dependence may arise in $`\varphi `$’s definition.
3. Calorons are BPS saturated (or self-dual) solutions to the Yang-Mills equations of motion in four-dimensional Euclidean spacetime (with time $`\tau `$ compactified on a circle). In particular, their energy-momentum tensor is precisely zero. The adjoint scalar field $`\varphi `$, being composed of noninteracting calorons, inherits this property, which in turn is expressed through a BPS equation for $`\varphi `$’s time dependence,
$$_\tau \varphi =V^{(1/2)}.$$
(3.2)
Here, $`V^{(1/2)}`$ is a ’square-root’ of the potential $`V(\varphi )=\mathrm{tr}(V^{(1/2)})^{}V^{(1/2)}`$, which governs the dynamics of $`\varphi `$. The potential $`V(\varphi )`$ is determined by the demand for BPS saturation.
In Sec. 3.2 we will give a microscopic definition for $`\varphi `$ and after evaluation see that this field is (up to a global gauge rotation) of the form
$$\varphi (\tau )=\sqrt{\frac{\mathrm{\Lambda }^3}{2\pi T}}\lambda _1\mathrm{exp}(2\pi iT\lambda _3\tau ).$$
(3.3)
The potential is found to be
$$V(\varphi )=\mathrm{\Lambda }^6\mathrm{tr}\varphi ^2=4\pi T\mathrm{\Lambda }^3.$$
(3.4)
Here, $`\mathrm{\Lambda }`$ is a fixed mass scale generated by dimensional transmutation. The modulus $`|\varphi |`$ falls off as $`\frac{1}{\sqrt{T}}`$. This dependence shows that the nontrivial-topology sector is strongly suppressed at large temperature.
Effective theory. Minimally coupling of $`\varphi `$ to the (up to now not included) topologically trivial fluctuations $`a_\mu `$ results in the effective action for the electric phase,
$$S_E=_0^{1/T}𝑑\tau d^3x\left(\frac{1}{2}\mathrm{tr}G_{\mu \nu }G_{\mu \nu }+\mathrm{tr}𝒟_\mu \varphi 𝒟_\mu \varphi +\mathrm{\Lambda }^6\varphi ^2\right),$$
(3.5)
where
$$G_{\mu \nu }^a=_\mu a_\nu ^a_\nu a_\mu ^ae\epsilon ^{abc}a_\mu ^ba_\nu ^c$$
(3.6)
is the field strength of topologically trivial fluctuations $`a_\mu `$,
$$𝒟_\mu \varphi =_\mu +ie[\varphi ,a_\mu ]$$
(3.7)
is the covariant derivative, and $`e`$ denotes the effective gauge coupling. There is no reason why the effective gauge coupling constant should be equal to the coupling constant of the fundamental theory.
The field $`\varphi `$ is seen to be quantum mechanically and statistically inert: The mass associated with its excitations is much larger than both temperature and the scale of admissible quantum fluctuations of $`|\varphi |`$ itself,
$$\frac{_{|\varphi |}^2V}{T^2}=12\pi ^2\text{and}\frac{_{|\varphi |}^2V}{|\varphi |^2}=3\lambda ^3,$$
(3.8)
where $`\lambda =\frac{2\pi T}{\mathrm{\Lambda }}`$ is the dimensionless temperature. In the electric phase of the theory, $`\lambda >11.65`$ (see below), so that $`3\lambda ^3`$ indeed is a large number. The quantum mechanically and thermodynamically stabilized field $`\varphi `$ can now be taken as a background for the equation of motion governing topologically trivial fluctuations,
$$𝒟_\mu G_{\mu \nu }=2ie[\varphi ,𝒟_\nu \varphi ].$$
(3.9)
This solution is to be part of the ground-state description and hence must not break the rotational invariance of the system. So it has to be a pure gauge. A pure-gauge solution to Eq. (3.9) reads
$$a_\mu ^{bg}=\delta _{\mu 4}\frac{\pi }{e}T\lambda _3.$$
(3.10)
It takes into account holonomy changing interactions between calorons mediated by topologically trivial fluctuations. Moreover, we have
$$𝒟_\mu \varphi =0$$
(3.11)
on the ground-state configuration $`\varphi `$, $`a_\mu ^{bg}`$. As a consequence, the action density in Eq. (3.5) evaluates to $`V(\varphi )`$ on the ground-state. This corresponds to a ground-state energy density $`\rho _{g.s.}=V(\varphi )=4\pi \mathrm{\Lambda }^3T`$ and a ground-state pressure $`P_{g.s.}=V(\varphi )=4\pi \mathrm{\Lambda }^3T`$.
The inert scalar $`\varphi `$ and the pure-gauge solution $`a_\mu ^{bg}`$ form the ground state about which loop expansions are performed. Therefore, the topologically trivial sector $`a_\mu `$ is split into the pure-gauge ground-state part $`a_\mu ^{bg}`$ and fluctuations $`\delta a_\mu `$ about this background. The adjoint scalar $`\varphi `$ renders some of the fluctuations $`\delta a_\mu `$ massive through the adjoint Higgs mechanism. In the case of SU(2), the gauge symmetry is dynamically broken<sup>1</sup><sup>1</sup>1 With the introduction of a composite field which breaks the gauge symmetry partially this approach conceptually resembles the Landau-Ginzburg-Abrikosov theory of superconductivity . to U(1). Two of the gauge bosons acquire the temperature dependent mass
$$m^2=4e^2|\varphi |^2=4e^2\frac{\mathrm{\Lambda }^3}{2\pi T},$$
(3.12)
while the third gauge boson is left massless<sup>2</sup><sup>2</sup>2 The calculation of the mass spectrum as well as other explicit calculations (such as the thermodynamical pressure in Chapter 4) are carried out after a transformation to unitary gauge, where $`a_\mu ^{bg}0`$ and hence
$$G_{\mu \nu }^a[a_\mu ]=G_{\mu \nu }^a[\delta a_\mu ].$$
The scalar field in this gauge takes the form
$$\varphi (\tau )=\sqrt{\frac{\mathrm{\Lambda }^3}{2\pi T}}\lambda _1.$$
It can be shown that the employed gauge transformation is admissible, see . . The former are referred to as tree-level heavy (TLH) modes, the latter as tree-level massless (TLM) modes. Since the masses of TLH excitations are temperature dependent, they are thermal quasiparticle fluctuations. The masses induced by the Higgs mechanism provide infrared cutoffs, so that no infrared divergences occur in loop expansions of thermodynamical quantities.
Compositeness constraints. The existence of the compositeness scale $`|\varphi |`$ present in the vacuum additionally implies cutoffs for the momenta propagating as vacuum fluctuations:
1. The offshellness of a quantum fluctuation must not exceed the scale $`|\varphi |`$,
$$|p^2m^2||\varphi |^2$$
(3.13)
where $`p`$ is the four-momentum of the quantum fluctuation and $`m`$ its mass.
2. The center-of-mass energy flowing into or out of a four-vertex must not exceed the scale $`|\varphi |`$,
$$\left|(p_1+p_2)^2\right|<|\varphi |^2,$$
(3.14)
where $`p_i`$ are the ingoing momenta.
The violation of any of these two conditions would immediately imply that the fluctuations generated out of the vacuum are able to destroy this vacuum, which is impossible. Both conditions Eqs. (3.13) and (3.14) have been formulated for a Minkowskian signature here, but can be rotated to the Euclidean.
By virtue of the cutoffs as they emerge above, the phase space for vacuum diagrams is strongly restricted. As a consequence of this, the interactions between quasiparticle fluctuations are very weak and (for sufficiently large temperatures) generate only tiny higher loop corrections to thermodynamical quantities. Moreover, ultraviolet divergences do not occur and the usual renormalization procedure is not needed in the effective theory.
Thermodynamical self-consistency. The thermodynamical quantities pressure, energy density, entropy density etc. as obtained from the fundamental SU(N) Yang-Mills theory are related by Legendre transforms. These relations hold in general, so their validity in the effective theory described by the action Eq. (3.5) has to be arranged for by imposing a condition of thermodynamical self-consistency. This condition has to assure that the $`T`$-derivatives of the TLH masses and the ground-state pressure cancel one another, so that only the explicit $`T`$-dependence arising from the Boltzmann weight enters $`T`$-derivatives of thermodynamical quantities. In particular, the relation
$$\rho =T\frac{dP}{dT}P$$
(3.15)
translates into
$$\frac{P}{a}=0,$$
(3.16)
where
$$a=2\pi e\lambda ^{3/2}=e\sqrt{\frac{\mathrm{\Lambda }^3}{2\pi T^3}}=\frac{m}{2T}$$
(3.17)
is a dimensionless measure for the quasiparticle mass. Eq. (3.16) results into an evolution equation $`\lambda (a)`$ for temperature as a function of the tree-level gauge boson mass. The evolution has two fixed points, namely $`a=0`$ and $`a=\mathrm{}`$. They signal the existence of both a lowest and a highest attainable temperature in the electric phase. These are denoted as $`\lambda _c=\lambda (a=\mathrm{})`$ and $`\lambda _P=\lambda (a=0)`$ respectively. The evolution is subject to the initial condition $`\lambda (a=0)\lambda _P`$. The low-temperature behavior of $`\lambda (a)`$ is practically independent of $`\lambda _P`$ as long as $`\lambda _P`$ is sufficiently large: The temperature at which the field $`\varphi `$ emerges does not influence the low-temperature physics. This is a signal of ultraviolet-infrared decoupling.
The evolution equation $`\lambda (a)`$ can be inverted to yield an evolution $`e(\lambda )`$ for the effective gauge coupling as a function of temperature. For SU(2) and SU(3) this is shown in Fig. 3.1: The critical temperature is $`\lambda _c=11.65`$ for SU(2), and the evolution exhibits a plateau, where the effective gauge coupling has the value $`e=5.1`$.
Phase transition to magnetic phase. At the critical temperature $`\lambda _c`$, the theory undergoes a second order phase transition from the electric to the so-called magnetic phase. The effective gauge coupling shows a logarithmic pole of the form
$$e(\lambda )\mathrm{log}(\lambda \lambda _c).$$
(3.18)
Hence the mass of the constituent BPS monopoles liberated by dissociating nontrivial holonomy calorons, given as
$$M_{\text{monopole}}\frac{4\pi }{e}T,$$
(3.19)
approaches zero with the consequence that magnetic monopoles condense.
In total, the theory is seen to have three phases: the electric phase at high temperatures, the magnetic phase for a small range of temperatures comparable to the scale $`\mathrm{\Lambda }`$, and a center phase for low temperatures. The electric phase is deconfining, the magnetic phase is preconfining, and the center phases completely confining. Here, we are only addressing the physics in the electric phase.
### 3.2 The composite adjoint Higgs field $`\varphi `$
In this section, the composite adjoint scalar field $`\varphi \varphi ^a\lambda ^a`$ ($`\lambda ^a`$ denote the Pauli matrices) which represents the topologically nontrivial, BPS saturated and trivial-holonomy part of the ground state is to be computed.
The field $`\varphi `$ can be written as a product of modulus and phase. As we are aiming for a description of a thermodynamical ground state, and the modulus $`|\varphi |`$ is a (gauge invariant and hence) physical quantity, it has to be homogenous both in space and time. To be able to calculate the modulus $`|\varphi |`$, the Yang-Mills scale must be known. The phase $`\frac{\varphi ^a}{|\varphi |}`$ is a dimensionless quantity, so that for its calculation no information about the Yang-Mills scale $`\mathrm{\Lambda }`$ is necessary. The classical caloron action $`S=\frac{8\pi ^2}{g^2}`$ is independent of temperature; this excludes explicit $`\beta `$-dependences of $`\varphi `$’s phase. The phase may depend on the temperature only via the periodicity of the caloron. Briefly, we can write
$$\varphi ^a=|\varphi |(\mathrm{\Lambda },\beta )\frac{\varphi ^a}{|\varphi |}\left(\frac{\tau }{\beta }\right).$$
(3.20)
The calorons generating the field $`\varphi `$ are BPS saturated; as an immediate consequence, the energy-momentum tensor vanishes identically on a caloron. As already mentioned, the field $`\varphi `$ will not include interactions, so the energy-momentum tensor has to vanish on $`\varphi `$ as well. This in turn can be described by a BPS equation for $`\varphi `$,
$$_\tau \varphi =V^{(1/2)},$$
(3.21)
where $`V^{(1/2)}`$ is the ’square-root’ of the potential $`V\mathrm{tr}(V^{(1/2)})^{}V^{(1/2)}`$.
In sections 3.2.1 and 3.2.3, we will define an adjointly transforming integral over a two-point function and demand that $`\varphi `$’s phase obeys the same equation of motion. In section 3.2.4 the BPS equation for $`\varphi `$ will be determined. In section 3.2.5 we will see that introducing the Yang-Mills scale $`\mathrm{\Lambda }`$ externally allows us to obtain $`\varphi `$’s modulus in terms of $`\mathrm{\Lambda }`$ and $`\beta `$. With phase and modulus known, the BPS equation is used to determine the potential governing the dynamics of $`\varphi `$. The potential is unique, if analytical dependence of the right-hand side of the BPS equation on the field $`\varphi `$ is demanded.
#### 3.2.1 Definition
We define the phase of the field $`\varphi `$ as proposed in by
$$\frac{\varphi ^a}{|\varphi |}\left(\frac{\tau }{\beta }\right)d^3x𝑑\rho \mathrm{tr}\lambda ^aF_{\mu \nu }(\tau ,0)\{(\tau ,0),(\tau ,𝒙)\}F_{\mu \nu }(\tau ,𝒙)\{(\tau ,𝒙),(\tau ,0)\},$$
(3.22)
where
$$\begin{array}{cc}\hfill |\varphi |^2& \frac{1}{2}\mathrm{tr}\varphi ^2\hfill \\ \hfill \{(\tau ,0),(\tau ,𝒙)\}& 𝒫\mathrm{exp}\left[i_{(\tau ,0)}^{(\tau ,𝒙)}𝑑z_\mu A_\mu (z)\right]\hfill \\ \hfill \{(\tau ,𝒙),(\tau ,0)\}& 𝒫\mathrm{exp}\left[i_{(\tau ,0)}^{(\tau ,𝒙)}𝑑z_\mu A_\mu (z)\right].\hfill \end{array}$$
(3.23)
The Wilson lines in Eq. (3.23) are to be calculated along the straight line connecting the points $`(\tau ,0)`$ and $`(\tau ,𝒙)`$. $`𝒫`$ denotes the path-ordering symbol, and $`\mathrm{tr}`$ the SU(2) trace. In (3.22), the dependences of the integrand on the caloron scale parameter $`\rho `$ and inverse temperature $`\beta `$ suppressed.
Eq. (3.22) is a definition of the phase $`\frac{\varphi ^a}{|\varphi |}`$ in the following sense: We demand that $`\frac{\varphi ^a}{|\varphi |}`$ obeys the same homogenous evolution equation in $`\tau `$ as the right hand side of (3.22) does,
$$𝒟\frac{\varphi }{|\varphi |}=0,$$
(3.24)
where $`𝒟`$ is a differential operator in $`\tau `$. Thus (3.24) is an equation of motion for $`\varphi `$’s phase.
The right-hand side of (3.22) is to be evaluated both on the caloron-field and the anticaloron-field and afterwards the sum is to be taken. We will see in the course of the calculation, that the definition (3.22) contains quite a number of ambiguities which span the solution space of the differential operator $`𝒟`$.
Under a gauge transformation $`\mathrm{\Omega }(\tau ,𝒙)`$, the involved objects transform as follows:
$$\begin{array}{cc}\hfill \{(\tau ,0),(\tau ,𝒙)\}& \mathrm{\Omega }^{}(\tau ,0)\{(\tau ,0),(\tau ,𝒙)\}\mathrm{\Omega }(\tau ,𝒙)\hfill \\ \hfill \{(\tau ,𝒙),(\tau ,0)\}& \mathrm{\Omega }^{}(\tau ,𝒙)\{(\tau ,𝒙),(\tau ,0)\}\mathrm{\Omega }(\tau ,0)\hfill \\ \hfill F_{\mu \nu }(\tau ,𝒙)& \mathrm{\Omega }^{}(\tau ,𝒙)F_{\mu \nu }(\tau ,𝒙)\mathrm{\Omega }(\tau ,𝒙).\hfill \end{array}$$
(3.25)
Hence the right-hand side of (3.22) indeed transforms like an adjoint scalar, namely
$$\frac{\varphi ^a}{|\varphi |}(\tau )R^{ab}(\tau )\frac{\varphi ^b}{|\varphi |}(\tau ),$$
(3.26)
where $`R^{ab}`$ is the SO(3) matrix
$$R^{ab}(\tau )\lambda ^b=\mathrm{\Omega }(\tau ,0)\lambda ^a\mathrm{\Omega }^{}(\tau ,0).$$
(3.27)
The field $`\varphi `$ only transforms under the time dependent part of the gauge transformation, the spatial dependence of the gauge transformation is lost in the macroscopic description.
The construction in (3.22) is invariant under spatial translations and hence the integration over spatial translations is trivial. For a gauge variant density as the integrand in (3.22) is, averaging over global color rotations yields zero and thus is forbidden. The same applies to time translations. The only modulus of the caloron solution that is integrated over is the scale parameter $`\rho `$ with flat measure.
#### 3.2.2 Are there alternative possibilities for the definition of $`\varphi `$’s phase?
The definition (3.22) for $`\varphi `$’s phase is not at all arbitrarily chosen. Indeed, trying to generalize the right-hand side of (3.22) always requires the introduction of either explicit temperature dependences (which are not allowed) or additional scales (which on the classical level do not exist):
* Every local definition including the field strength $`F_{\mu \nu }`$ only, such as
$$\mathrm{tr}\lambda ^aF_{\mu \nu }F_{\nu \kappa }F_{\kappa \mu },\mathrm{tr}\lambda ^aF_{\mu \nu }F_{\nu \kappa }F_{\kappa \rho }F_{\rho \mu }\text{etc.}$$
(3.28)
or
$$\mathrm{tr}\epsilon _{abc}F_{\mu \nu }\lambda ^bF_{\nu \kappa }\lambda ^cF_{\kappa \mu }\text{etc.}$$
(3.29)
yields zero when evaluated on the (anti)caloron field. This is due to the (anti)self-duality of the (anti)caloron field Eq. (2.98).
* One could be tempted to consider higher $`n`$-point functions of the type employed in Eq. (3.22), such as
$$\begin{array}{c}\hfill \beta ^1d^3xd^3y𝑑\rho \mathrm{tr}\lambda ^aF_{\mu \nu }(\tau ,0)\{(\tau ,0),(\tau ,𝒙)\}F_{\nu \kappa }(\tau ,𝒙)\\ \hfill \{(\tau ,𝒙),(\tau ,𝒚)\}F_{\kappa \mu }(\tau ,𝒚)\{(\tau ,𝒚),(\tau ,0)\}\end{array}$$
(3.30)
In such an $`n`$-point function, one has to introduce an additional factor $`\beta ^{2n}`$ to get a dimensionless object. $`\varphi `$’s phase is supposed to depend on temperature only via the temperature dependence of the caloron field, but not explicitly. (This is due to the temperature-independence of the caloron action.) Therefore the $`n`$-point function (3.30) and its generalization for $`n3`$ are forbidden.
* One could think of shifting the spatial part of the starting point of the construction in (3.22) from 0 to an arbitrary point $`𝒛`$. Any given value of $`|𝒛|0`$ would introduce an additional scale; but, on the classical level, such a scale does not exist. The same argument applies to replacing the straight Wilson lines with curved arcs. Their curvature again corresponds to an additional scale which is physically not present.
* The definition (3.22) does not include caloron solutions with topological charge $`|Q|>1`$. This has the following reason: A caloron of topological charge $`|Q|>1`$ has $`m>1`$ dimensionful moduli. For example, a caloron with $`|Q|=2`$ has three dimensionful moduli, namely two scale parameters and the distance between its two centers. An $`n`$-point function of the type displayed in Eq. (3.30) contains $`n`$ field strength tensors (mass dimension 2) and $`n1`$ integrations over 3-space (mass dimension $`3`$). The $`m`$ dimensionful moduli of the caloron have to be integrated as well (mass dimension $`1`$ each). These combine to mass dimension $`2n3(n1)m=3nm`$. We are looking for a dimensionless object without any explicit dependence on $`\beta `$. But this is not possible with $`n>2`$ and $`m>1`$. Therefore, calorons of higher topological charge are excluded.
#### 3.2.3 Calculation
In this section, we want to evaluate the right-hand side of (3.22),
$$d^3x𝑑\rho \mathrm{tr}\lambda ^aF_{\mu \nu }(\tau ,0)\{(\tau ,0),(\tau ,𝒙)\}F_{\mu \nu }(\tau ,𝒙)\{(\tau ,𝒙),(\tau ,0)\},$$
(3.31)
on the single caloron solution (cf. Sec. 2)
$$A_\mu (\tau ,𝒙)=\overline{\eta }_{\mu \nu }^a\frac{\lambda ^a}{2}_\nu \mathrm{ln}\mathrm{\Pi }(\tau ,r)$$
(3.32)
with the pre-potential $`\mathrm{\Pi }(\tau ,r)`$ given as
$$\mathrm{\Pi }(\tau ,r)=1+\frac{\pi \rho ^2}{\beta r}\frac{\mathrm{sinh}\frac{2\pi r}{\beta }}{\mathrm{cosh}\frac{2\pi r}{\beta }\mathrm{cos}\frac{2\pi \tau }{\beta }},$$
(3.33)
and $`r|𝒙|`$.
The Wilson lines. The single caloron solution has a hedgehog like behavior in the sense that the spatial part of the scalar product $`x_\mu A_\mu (\tau ,𝒙)`$ has the same orientation both in 3-space and color-space,
$$\begin{array}{cc}\hfill x_iA_i^a(\tau ,𝒙)& =x_i\overline{\eta }_{i\nu }^a_\nu \mathrm{ln}\mathrm{\Pi }(\tau ,r)\hfill \\ & =x_i\left(\epsilon _{ai\nu }\delta _{\nu 4}\delta _{ai}\right)_\nu \mathrm{ln}\mathrm{\Pi }(\tau ,r)\hfill \\ & =\epsilon _{ain}\frac{x_ix_n}{r}_r\mathrm{ln}\mathrm{\Pi }(\tau ,r)+x_i\delta _{ai}_4\mathrm{ln}\mathrm{\Pi }(\tau ,r)\hfill \\ & =x_a_4\mathrm{ln}\mathrm{\Pi }(\tau ,r).\hfill \end{array}$$
(3.34)
This property allows to discard the path ordering operation in the calculation of the Wilson line $`\{(\tau ,0),(\tau ,𝒙)\}`$ defined in Eq. (3.23). We parameterize the path as $`z_\mu (s)=(\tau ,s𝒙)`$ with $`0s1`$. Hence we have
$$\begin{array}{cc}\hfill \{(\tau ,0),(\tau ,𝒙)\}& =𝒫\mathrm{exp}\left[i_{(\tau ,0)}^{(\tau ,𝒙)}𝑑z_\mu A_\mu (z)\right]\hfill \\ & =𝒫\mathrm{exp}\left[i_0^1𝑑sx_iA_i(\tau ,s𝒙)\right]\hfill \\ & =𝒫\mathrm{exp}\left[i_0^1𝑑s\frac{x_i}{2}\lambda _i_4\mathrm{ln}\mathrm{\Pi }(\tau ,sr)\right]\hfill \\ & =\mathrm{exp}\left[i\lambda ^i\frac{x_i}{2}_0^1𝑑s_4\mathrm{ln}\mathrm{\Pi }(\tau ,sr)\right]\hfill \\ & =\mathrm{cos}\left(g(\tau ,r)\right)+i\lambda ^i\frac{x_i}{r}\mathrm{sin}\left(g(\tau ,r)\right),\hfill \end{array}$$
(3.35)
where we define
$$g(\tau ,r)_0^1𝑑s\frac{r}{2}_4\mathrm{ln}\mathrm{\Pi }(\tau ,sr).$$
(3.36)
The evaluation of the Wilson line on the single anticaloron solution yields just a change of sign in the argument of the exponential function. So, for caloron and anticaloron we have
$$\{(\tau ,0),(\tau ,𝒙)\}_{C,A}=\mathrm{cos}\left(g(\tau ,r)\right)\pm i\lambda _i\frac{x_i}{r}\mathrm{sin}\left(g(\tau ,r)\right)$$
(3.37)
respectively. The following relations hold:
$$\begin{array}{cc}\hfill \{(\tau ,0),(\tau ,𝒙)\}_C& =\{(\tau ,𝒙),(\tau ,0)\}_A=\{(\tau ,0),(\tau ,𝒙)\}_A=\{(\tau ,0),(\tau ,𝒙)\}_C^{}\hfill \\ & =\{(\tau ,𝒙),(\tau ,0)\}_C^{}=\{(\tau ,0),(\tau ,𝒙)\}_A^{}=\{(\tau ,𝒙),(\tau ,0)\}_C.\hfill \end{array}$$
(3.38)
The integrand of Eq. (3.36) for large $`r`$ behaves like a $`\delta `$-function in $`s`$,
$$\underset{r\mathrm{}}{lim}\frac{r}{2}_4\mathrm{ln}\mathrm{\Pi }(\tau ,sr)=\delta (s)f(\tau ).$$
(3.39)
This property has been established numerically.
Caloron field strength. The field strength $`F_{\mu \nu }`$ on a caloron can be calculated as
$`F_{\mu \nu }^a`$ $`=_\mu A_\nu ^a_\nu A_\mu ^ai[A_\mu ,A_\nu ]^a`$ (3.40)
$`=\overline{\eta }_{\mu \nu }^a{\displaystyle \frac{(_\kappa \mathrm{\Pi })(_\kappa \mathrm{\Pi })}{\mathrm{\Pi }^2}}\overline{\eta }_{\nu \kappa }^a{\displaystyle \frac{\mathrm{\Pi }(_\mu _\kappa \mathrm{\Pi })2(_\mu \mathrm{\Pi })(_\kappa \mathrm{\Pi })}{\mathrm{\Pi }^2}}+\overline{\eta }_{\mu \kappa }^a{\displaystyle \frac{\mathrm{\Pi }(_\nu _\kappa \mathrm{\Pi })2(_\nu \mathrm{\Pi })(_\kappa \mathrm{\Pi })}{\mathrm{\Pi }^2}}.`$
The calculation is performed in Appendix A.2. To obtain the field strength of an anticaloron, $`\overline{\eta }`$ has to be replaced with $`\eta `$. In particular, the field strength at $`𝒙=0`$ is
$$F_{\mu \nu }^a(\tau ,0)=\eta _{\mu \nu }^a\left(\frac{[_4\mathrm{\Pi }(\tau ,0)]^2}{\mathrm{\Pi }^2(\tau ,0)}\frac{2}{3}\frac{_4^2\mathrm{\Pi }(\tau ,0)}{\mathrm{\Pi }(\tau ,0)}\right).$$
(3.41)
The integrand of Eq. (3.31). Inserting the result for the Wilson lines Eq. (3.37) into the expression (3.31) and writing the field strength in components, $`F_{\mu \nu }=F_{\mu \nu }^a\frac{\lambda ^a}{2}`$, we see that
$$\begin{array}{cc}& \mathrm{tr}\lambda ^aF_{\mu \nu }(\tau ,0)\{(\tau ,0),(\tau ,𝒙)\}F_{\mu \nu }(\tau ,𝒙)\{(\tau ,𝒙),(\tau ,0)\}|_{\text{Caloron}}\hfill \\ & =\frac{1}{2}\mathrm{tr}\left[\lambda ^a\lambda ^b\left(\mathrm{cos}g(\tau ,r)+i\lambda ^c\frac{x^c}{r}\mathrm{sin}g(\tau ,r)\right)\lambda ^d\left(\mathrm{cos}g(\tau ,r)i\lambda ^e\frac{x^e}{r}\mathrm{sin}g(\tau ,r)\right)\right]\hfill \\ & F_{\mu \nu }^b(\tau ,0)F_{\mu \nu }^d(\tau ,𝒙)\hfill \end{array}$$
(3.42)
is to be computed. Performing the trace and contracting Lorentz and color indices is a straightforward but somewhat lengthy calculation; for details see Appendix A.3. The result is<sup>3</sup><sup>3</sup>3 Note that $`_r`$ denotes the derivative with respect to the radial coordinate $`r|𝒙|`$. For more notations and conventions, see Appendix A.1.
$$\begin{array}{cc}& \mathrm{tr}\lambda ^aF_{\mu \nu }(\tau ,0)\{(\tau ,0),(\tau ,𝒙)\}F_{\mu \nu }(\tau ,𝒙)\{(\tau ,𝒙),(\tau ,0)\}|_{\text{Caloron}}\hfill \\ & =2i\frac{x^a}{r}\left(\frac{[_4\mathrm{\Pi }(\tau ,0)]^2}{\mathrm{\Pi }^2(\tau ,0)}\frac{2}{3}\frac{_4^2\mathrm{\Pi }(\tau ,0)}{\mathrm{\Pi }(\tau ,0)}\right)\hfill \\ & \{2\mathrm{cos}(2g(\tau ,r))(2\frac{[_4\mathrm{\Pi }(\tau ,r)][_r\mathrm{\Pi }(\tau ,r)]}{\mathrm{\Pi }^2(\tau ,r)}\frac{_4_r\mathrm{\Pi }(\tau ,r)}{\mathrm{\Pi }(\tau ,r)})\hfill \\ & +\mathrm{sin}(2g(\tau ,r))(2\frac{[_r\mathrm{\Pi }(\tau ,r)]^2}{\mathrm{\Pi }^2(\tau ,r)}2\frac{[_4\mathrm{\Pi }(\tau ,r)]^2}{\mathrm{\Pi }^2(\tau ,r)}+\frac{_4^2\mathrm{\Pi }(\tau ,r)}{\mathrm{\Pi }(\tau ,r)}\frac{_r^2\mathrm{\Pi }(\tau ,r)}{\mathrm{\Pi }(\tau ,r)})\}.\hfill \end{array}$$
(3.43)
The factor
$$\frac{[_4\mathrm{\Pi }(\tau ,0)]^2}{\mathrm{\Pi }^2(\tau ,0)}\frac{2}{3}\frac{_4^2\mathrm{\Pi }(\tau ,0)}{\mathrm{\Pi }(\tau ,0)}=\frac{16\pi ^4}{3}\frac{\rho ^2}{\beta ^2}\frac{\pi ^2\rho ^2+\beta ^2\left(2+\mathrm{cos}\frac{2\pi \tau }{\beta }\right)}{\left(2\pi ^2\rho ^2+\beta ^2\left(1\mathrm{cos}\frac{2\pi \tau }{\beta }\right)\right)^2}$$
(3.44)
arises from the field strength at the point $`(\tau ,0)`$, cf. Eq. (3.41), and contains no dependence on $`r`$. Note that the expression (3.43) is proportional to the (spatial) unit vector $`\frac{x^a}{r}`$.
Integration over position space. The integration over position space demanded in (3.31) is performed in polar coordinates. On the one hand, the angular integration over the unit vector $`\frac{x^a}{r}`$ yields zero. But on the other hand, the radial integral is infinite. This can be seen as follows: The pre-potential (3.33) for large $`r`$ behaves as
$$\mathrm{\Pi }(\tau ,r)1+\frac{\pi \rho ^2}{\beta r}\text{(}r\beta /2\pi \text{)},$$
(3.45)
and its second spatial derivative approaches
$$_r^2\mathrm{\Pi }(\tau ,r)_r^2\left(1+\frac{\pi \rho ^2}{\beta r}\right)=\frac{2\pi \rho ^2}{\beta r^3}\text{(}r\beta /2\pi \text{)}.$$
(3.46)
Hence the very last term in Eq. (3.43) has the asymptotic behavior
$$\frac{_r^2\mathrm{\Pi }(\tau ,r)}{\mathrm{\Pi }(\tau ,r)}\frac{2\pi \rho ^2}{\beta r^3}\text{(}r\beta /2\pi \text{)}$$
(3.47)
and gives rise to the logarithmically divergent integral<sup>4</sup><sup>4</sup>4 Note that the integral is convergent at $`r=0`$, since we have
$$\underset{r0}{lim}\frac{_r^2\mathrm{\Pi }(\tau ,r)}{\mathrm{\Pi }(\tau ,r)}=\frac{4\pi ^4\rho ^2}{3\beta ^2}\frac{1}{\mathrm{sin}^2\left(\frac{\pi \tau }{\beta }\right)}\frac{\left(2+\mathrm{cos}\left(\frac{2\pi \tau }{\beta }\right)\right)}{2\pi ^2\rho ^2+\beta ^2\left(1\mathrm{cos}\left(\frac{2\pi \tau }{\beta }\right)\right)}.$$
$$\begin{array}{cc}\hfill _0^{\mathrm{}}𝑑rr^2\mathrm{sin}(2g(\tau ,r))\frac{_r^2\mathrm{\Pi }(\tau ,r)}{\mathrm{\Pi }(\tau ,r)}& =\text{finite}+_R^{\mathrm{}}𝑑r\mathrm{sin}(2g(\tau ,r))\frac{2\pi }{\beta r}\hfill \\ & =\text{finite}+\frac{2\pi \rho ^2}{\beta }\left(\underset{r\mathrm{}}{lim}\mathrm{sin}(2g(\tau ,r))\right)_R^{\mathrm{}}\frac{dr}{r},\hfill \end{array}$$
(3.48)
where $`R\beta /2\pi `$. The function $`g(r)`$ defined in Eq. (3.36) has a finite limit for $`r\mathrm{}`$, which is reached very rapidly. It has been taken out of the integral.
All the other terms in Eq. (3.43) give rise to finite integrals in $`r`$, and hence do not contribute: For the square of the first spatial derivative, we have
$$\frac{[_r\mathrm{\Pi }(\tau ,r)]^2}{\mathrm{\Pi }^2(\tau ,r)}\frac{\pi ^2\rho ^4}{\beta ^2r^4}\text{(}r\beta /2\pi \text{)},$$
(3.49)
which is convergent. All the other terms contain at least one time derivative and thus vanish exponentially, e. g.
$$_\tau \mathrm{\Pi }(\tau ,r)\frac{2\pi \rho ^2}{\beta r}\mathrm{sin}\left(\frac{2\pi \tau }{\beta }\right)\mathrm{exp}\left(\frac{2\pi r}{\beta }\right)\text{(}r\beta /2\pi \text{)}.$$
(3.50)
Thus the only nonvanishing contribution to the expression (3.31) is
$`2i{\displaystyle 𝑑\rho \left(\frac{[_4\mathrm{\Pi }(\tau ,0)]^2}{\mathrm{\Pi }^2(\tau ,0)}\frac{2}{3}\frac{_4^2\mathrm{\Pi }(\tau ,0)}{\mathrm{\Pi }(\tau ,0)}\right)_{S_2}𝑑\mathrm{\Omega }\frac{x^a}{r}_0^{\mathrm{}}𝑑rr^2\mathrm{sin}(2g(\tau ,r))\frac{_r^2\mathrm{\Pi }(\tau ,r)}{\mathrm{\Pi }(\tau ,r)}}`$
$`={\displaystyle \frac{64i\pi ^5}{3}}{\displaystyle 𝑑\rho \frac{\rho ^4}{\beta ^3}\frac{\pi ^2\rho ^2+\beta ^2\left(2+\mathrm{cos}\frac{2\pi \tau }{\beta }\right)}{\left(2\pi ^2\rho ^2+\beta ^2\left(1\mathrm{cos}\frac{2\pi \tau }{\beta }\right)\right)^2}_{S_2}𝑑\mathrm{\Omega }\frac{x^a}{r}_R^{\mathrm{}}\frac{dr}{r}\mathrm{sin}(2g(\tau ,r))}.`$ (3.51)
Regularization. As already pointed out, this expression contains a product zero $`\times `$ infinity and thus needs to be given a finite value by prescribing a regularization procedure. The radial integration in Eq. (3.2.3) is regularized according to
$$_R^{\mathrm{}}\frac{dr}{r}\beta ^\epsilon _R^{\mathrm{}}\frac{dr}{r^{1+\epsilon }}$$
(3.52)
with $`\epsilon >0`$. This integral is
$$\begin{array}{cc}\hfill \beta ^\epsilon _R^{\mathrm{}}\frac{dr}{r^{1+\epsilon }}& =\beta ^\epsilon _0^{\mathrm{}}\frac{dr}{(r+R)^{1+\epsilon }}=\frac{1}{\epsilon }\left(\frac{\beta }{R}\right)^\epsilon \hfill \\ & =\frac{1}{\epsilon }\mathrm{log}\left(\frac{R}{\beta }\right)+\frac{1}{2}\epsilon \mathrm{log}^2\left(\frac{R}{\beta }\right)+𝒪(\epsilon ^2).\hfill \end{array}$$
(3.53)
The right-hand side is a regular expression for $`\epsilon 0`$. It can be regarded as the analytical continuation of the integral on the left-hand side for $`|\epsilon |1`$. When smearing the regulator $`\epsilon `$ over a small interval $`[\eta ,\eta ]`$ with $`0<\eta 1`$ as
$$\begin{array}{cc}& \frac{1}{2\eta }_\eta ^\eta 𝑑\epsilon \left(\frac{1}{\epsilon \pm i0}\mathrm{log}\left(\frac{R}{\beta }\right)+\frac{1}{2}\epsilon \mathrm{log}^2\left(\frac{R}{\beta }\right)+𝒪(\epsilon ^2)\right)=\frac{\pi i}{2\eta }\mathrm{log}\left(\frac{R}{\beta }\right)+𝒪(\eta ^2),\hfill \end{array}$$
(3.54)
an ambiguity appears: There are two possibilities how to circumvent the pole.
The angular integration is regularized via the introduction of a defect (or surplus) angle in the azimuthal integration,
$$_0^\pi 𝑑\theta \mathrm{sin}\theta _0^{2\pi }𝑑\phi _0^\pi 𝑑\theta \mathrm{sin}\theta _{\alpha _C\pm \eta ^{}}^{\alpha _C+2\pi \eta ^{}}𝑑\phi $$
(3.55)
with $`0<\eta ^{}1`$ and an arbitrary angle $`0\alpha _C<2\pi `$. The regulated radial integral yields
$$\begin{array}{cc}\hfill _0^\pi 𝑑\theta \mathrm{sin}\theta _{\alpha _C\pm \eta ^{}}^{\alpha _C+2\pi \eta ^{}}𝑑\phi \frac{x^a}{r}& =\pi \mathrm{sin}\eta ^{}(\delta _{a1}\mathrm{cos}\alpha _C+\delta _{a2}\mathrm{sin}\alpha _C)\hfill \\ & =\pi \eta ^{}(\delta _{a1}\mathrm{cos}\alpha _C+\delta _{a2}\mathrm{sin}\alpha _C)+𝒪(\eta ^2).\hfill \end{array}$$
(3.56)
The regularization prescription Eq. (3.55) clearly singles out the axis with unit vector $`(\mathrm{cos}\alpha _C,\mathrm{sin}\alpha _C,0)`$. Later this will be seen to relate to a global gauge choice.
At the end of the calculation, both regularization parameters $`\eta `$ and $`\eta ^{}`$ have to be sent to zero. Their ratio in this limit will be denoted by $`\mathrm{\Xi }_C`$,
$$\underset{\eta ,\eta ^{}0}{lim}\frac{\eta ^{}}{\eta }=\mathrm{\Xi }_C.$$
(3.57)
This is a positive but otherwise unknown number. Inserting the regularized expressions Eqs. (3.54) and (3.56) and into (3.2.3), we get
$$\begin{array}{cc}& d^3x𝑑\rho \mathrm{tr}\lambda ^aF_{\mu \nu }(\tau ,0)\{(\tau ,0),(\tau ,𝒙)\}F_{\mu \nu }(\tau ,𝒙)\{(\tau ,𝒙),(\tau ,0)\}|_{\text{Caloron}}\hfill \\ & =\pm \mathrm{\Xi }_C(\delta _{a1}\mathrm{cos}\alpha _C+\delta _{a2}\mathrm{sin}\alpha _C)𝒜\left(\frac{2\pi \tau }{\beta }\right)\hfill \end{array}$$
(3.58)
where the dimensionless function $`𝒜`$ is defined as
$$𝒜\left(\frac{2\pi \tau }{\beta }\right)\frac{32}{3}\frac{\pi ^7}{\beta ^3}𝑑\rho \left(\underset{r\mathrm{}}{lim}\mathrm{sin}2g(\tau ,r)\right)\rho ^4\frac{\pi ^2\rho ^2+\beta ^2\left(2+\mathrm{cos}\frac{2\pi \tau }{\beta }\right)}{\left(2\pi ^2\rho ^2+\beta ^2\left(1\mathrm{cos}\frac{2\pi \tau }{\beta }\right)\right)^2}.$$
(3.59)
Integration over scale parameter. Up to now, we did not specify the range for the integration over the scale parameter $`\rho `$. As the integrand in Eq. (3.59) asymptotically behaves like $`\rho ^2`$, integrating $`\rho `$ from zero to infinity yields an infinite expression. To see what is going on, a cutoff is introduced in units of inverse temperature $`\beta `$,
$$𝑑\rho _0^{\zeta \beta }𝑑\rho (\zeta >0).$$
(3.60)
This generates an additional dependence of $`𝒜`$ on $`\zeta `$. For $`\zeta 1`$, we will have $`𝒜\zeta ^3`$. The integral in Eq. (3.59) has to be evaluated numerically. Due to the property Eq. (3.39), the limit $`lim_r\mathrm{}\mathrm{sin}2g(\tau ,r)`$ is reached very fast; for our purposes, putting $`r=10`$ in numerical calculations is fully sufficient.
Fig. 3.2 shows the $`\tau `$-dependence of $`𝒜`$ for various values of $`\zeta `$. With growing $`\zeta `$, the function $`𝒜`$ rapidly approaches a sine curve,
$$𝒜(\frac{2\pi \tau }{\beta },\zeta )272\zeta ^3\mathrm{sin}\left(\frac{2\pi \tau }{\beta }\right)\text{(}\zeta \mathrm{}\text{)}.$$
(3.61)
Already for $`\zeta =10`$, the difference between the two curves could not be resolved any more in the figure. The prefactor 272 has been fitted numerically, see Table 3.1.
Therefore, the evaluation of the expression (3.31) on the single caloron yields
$$\begin{array}{cc}& d^3x𝑑\rho \mathrm{tr}\lambda ^aF_{\mu \nu }(\tau ,0)\{(\tau ,0),(\tau ,𝒙)\}F_{\mu \nu }(\tau ,𝒙)\{(\tau ,𝒙),(\tau ,0)\}|_{\text{Caloron}}\hfill \\ & =\pm 272\zeta ^3\mathrm{\Xi }_C(\delta _{a1}\mathrm{cos}\alpha _C+\delta _{a2}\mathrm{sin}\alpha _C)\mathrm{sin}\left(\frac{2\pi \tau }{\beta }\right).\hfill \end{array}$$
(3.62)
#### 3.2.4 Ambiguities and BPS saturation
The expression (3.31) has been evaluated on the caloron in the previous section. Evaluation on the anticaloron yields the same result, if we agree upon circumventing the pole appearing in Eq. (3.54) in the opposite way as compared to the caloron. According to the definition, both contributions are to be added. In this process some ambiguities occur:
1. The undefined number $`lim_{\eta ,\eta ^{}0}\frac{\eta ^{}}{\eta }`$ needs not necessarily be the same in the caloron and anticaloron case, we chose to call it $`\mathrm{\Xi }_A`$ for the anticaloron (instead of $`\mathrm{\Xi }_C`$ for the caloron).
2. The same applies to the arbitrarily chosen axis singled out by angular regularization. It is no restriction of generality to assume that both axes lie in the $`x_1x_2`$-plane, but with different azimuthal angles $`\alpha _C`$ and $`\alpha _A`$.
3. As (3.22) is supposed to define only the equation of motion for $`\varphi `$, we may as well introduce a shift in time; this does not change the operator $`𝒟`$. Again, this shift may be different for the two contributions, namely $`\tau \tau +\tau _C`$ for the caloron and $`\tau \tau +\tau _A`$ for the anticaloron.
The requested sum of caloron and anticaloron contribution (including all the above ambiguities) is then
$$\begin{array}{cc}\hfill \frac{\varphi ^a}{|\varphi |}272\zeta ^3\{& \pm \mathrm{\Xi }_C(\delta _{a1}\mathrm{cos}\alpha _C+\delta _{a2}\mathrm{sin}\alpha _C)\mathrm{sin}\left(\frac{2\pi (\tau +\tau _C)}{\beta }\right)\hfill \\ & \pm \mathrm{\Xi }_A(\delta _{a1}\mathrm{cos}\alpha _A+\delta _{a2}\mathrm{sin}\alpha _A)\mathrm{sin}\left(\frac{2\pi (\tau +\tau _A)}{\beta }\right)\}.\hfill \end{array}$$
(3.63)
The caloron (or anticaloron) contribution in Eq. (3.63) alone represents a linearly polarized harmonic oscillation in adjoint color space. The three ambiguities given above can be viewed as the free parameters of such an oscillation, namely modulus, phase-shift and polarization axis. The sum of caloron and anticaloron contribution Eq. (3.63) taking into account all the above ambiguities is an elliptically polarized oscillation in adjoint color space. Here, the polarization plane is the $`x_1x_2`$-plane; this is only due to our choice of the angular regularization and no physical property.
The right-hand side of Eq. (3.63) obviously is annihilated by the second order differential operator
$$𝒟=_\tau ^2+\left(\frac{2\pi }{\beta }\right)^2.$$
(3.64)
Note that the ambiguities inherent to the definition (3.22) span the solution space of $`𝒟`$.
Imposing BPS saturation. The field $`\varphi `$ is composed of noninteracting trivial-holonomy calorons; the caloron itself is an energy- and pressure-free, BPS saturated configuration and no interaction whatsoever has been included. Hence the composite object must again be energy- and pressure-free and BPS saturated; that means it does not only obey a second order differential equation (the equation of motion), but also a first order differential equation (the BPS equation). Thus we need to find first-order equations whose solutions solve the second order equation
$$_\tau ^2\varphi +\left(\frac{2\pi }{\beta }\right)^2\varphi =0$$
(3.65)
as well. There are two such equations<sup>5</sup><sup>5</sup>5 The solutions to the equations $`_\tau \varphi =\pm \frac{2\pi i}{\beta }\varphi `$ as well solve Eq. (3.65); in spite of this they are not allowed here because they are not adjoint fields., namely
$$_\tau \varphi =\pm \frac{2\pi i}{\beta }\lambda _3\varphi .$$
(3.66)
Choosing any (normalized) linear combination of Pauli matrices instead of $`\lambda _3`$ would be possible equally well: Eq. (3.66) is subject to a global gauge ambiguity. The solutions to (3.66) are given as
$$\varphi =C\lambda _1\mathrm{exp}\left(\pm \frac{2\pi i}{\beta }\lambda _3(\tau \tau _0)\right),$$
(3.67)
where $`C`$ and $`\tau _0`$ are real integration constants. This solution is a circularly polarized oscillation. It winds along an $`S_1`$ in the group manifold $`S_3`$ of SU(2).
Thus the demand for BPS saturation forces an elliptical polarization in Eq. (3.63) into a circular polarization. The undetermined and formerly independent quantities are now subject to the relations
$`\mathrm{\Xi }_C`$ $`=\mathrm{\Xi }_A,`$ $`\tau _C\tau _A`$ $`=\pm {\displaystyle \frac{\pi }{2}},`$ $`\alpha _C\alpha _A`$ $`=\pm {\displaystyle \frac{\pi }{2}}.`$ (3.68)
The modulus of the oscillation and its phase-shift are still undetermined constants of integration. The former will be considered in the following section, the latter is of no physical significance. In addition, a global gauge ambiguity, i. e. the plane in which the oscillation takes place, is still present.
#### 3.2.5 Obtaining $`\varphi `$’s modulus
Let us now assume the existence of an externally given scale $`\mathrm{\Lambda }`$ which determines $`\varphi `$’s modulus. We allow for explicit dependence of $`\varphi `$ on the scale $`\mathrm{\Lambda }`$, the temperature $`\beta `$ and on Euclidean time through $`\frac{\tau }{\beta }`$,
$$\varphi =\varphi (\beta ,\mathrm{\Lambda },\frac{\tau }{\beta }).$$
(3.69)
As the phase found in Eq. (3.67) shall be preserved even in case of the presence of a scale $`\mathrm{\Lambda }`$, the right-hand side of the BPS equation
$$_\tau \varphi =V^{(1/2)}$$
(3.70)
may only depend linearly on $`\varphi `$. Besides that, we demand an analytical dependence of $`V`$ on $`\varphi `$. The potential $`V`$ (and its ’square-root’ $`V^{(1/2)}`$) may depend on the temperature through the periodicity of $`\varphi `$. An explicit dependence on $`\beta `$ is not possible because no explicit $`\beta `$-dependence occurs in the average over the caloron moduli space.
These conditions leave only the two possibilities
$$_\tau \varphi =\pm i\mathrm{\Lambda }\lambda _3\varphi $$
(3.71)
and<sup>6</sup><sup>6</sup>6 Note that $`\frac{\varphi }{|\varphi |^2}=\varphi ^1=\varphi _0^1_{n=0}^{\mathrm{}}(1)^n\varphi _0^n(\varphi \varphi _0)^n`$ is indeed an analytical function of $`\varphi `$.
$$_\tau \varphi =\pm i\mathrm{\Lambda }^3\lambda _3\frac{\varphi }{|\varphi |^2}.$$
(3.72)
Using Eqs. (3.20) and (3.67), we write
$$\varphi =|\varphi (\beta ,\mathrm{\Lambda })|\lambda _1\mathrm{exp}\left(\pm \frac{2\pi i}{\beta }\lambda _3\tau \right).$$
(3.73)
The modulus of $`\varphi `$ is gauge invariant and hence a physical quantity describing a thermodynamical ground state, so it has to be homogenous in space and time. Inserting the decomposition (3.73) into the first BPS equation (3.71), we get
$$\mathrm{\Lambda }=\frac{2\pi }{\beta }.$$
(3.74)
This obviously can not be satisfied, since the scale $`\mathrm{\Lambda }`$ is a constant and $`\beta `$ is the inverse temperature. From the second possibility Eq. (3.72), we get
$$|\varphi |(\beta ,\mathrm{\Lambda })=\sqrt{\frac{\beta \mathrm{\Lambda }^3}{2\pi }}=\sqrt{\frac{\mathrm{\Lambda }^3}{2\pi T}},$$
(3.75)
which is no contradiction. So, Eq. (3.72) is the only acceptable BPS equation, and $`\varphi `$ has the form
$$\varphi =\sqrt{\frac{\beta \mathrm{\Lambda }^3}{2\pi }}\lambda _1\mathrm{exp}\left(\pm \frac{2\pi i}{\beta }\lambda _3\tau \right).$$
(3.76)
Eq. (3.75) shows that the field $`\varphi `$ is power suppressed in $`T`$, and hence all topologically nontrivial effects die off at high temperature. The right-hand side of the BPS equation defines the ’square-root’ of $`\varphi `$’s potential,
$$V(\varphi )=\mathrm{tr}(V^{(1/2)})^{}V^{(1/2)}=\mathrm{\Lambda }^6\mathrm{tr}\varphi ^2.$$
(3.77)
Under the above assumptions, the potential is unique. The Lagrangian for the field $`\varphi `$ is
$$=\mathrm{tr}(_\tau \varphi )^2+V(\varphi ).$$
(3.78)
The consequences of minimally coupling $`\varphi `$ to the topologically trivial sector have been investigated in ; some of the results are briefly reviewed in Sec. 3.1.
## Chapter 4 Thermodynamical Pressure in the Electric Phase
In this chapter the two-loop corrections to the thermodynamical pressure in the electric phase of SU(2) Yang-Mills theory are computed. In view of the evolution of the effective gauge coupling in the electric phase (cf. Fig. 3.1), one can constrain oneself to the case $`e>\frac{1}{2}`$. This will simplify the calculation. For a numerical evaluation, the plateau value $`e=5.1`$ is used. We will set up the prerequisites for the calculation in Sec. 4.1. The calculation is performed in Sec. 4.2. In Sec. 4.3, the two-loop contributions are compared to the one-loop result for the pressure.
### 4.1 Feynman rules and other prerequisites
#### 4.1.1 One-loop pressure
The thermodynamical pressure $`P`$ is defined as the derivative of the partition function $`Z`$ with respect to the volume of the system,
$$P=T\frac{\mathrm{ln}Z}{V}.$$
(4.1)
In a field theory, $`P`$ can be calculated order by order in a loop expansion with the help of diagrammatic techniques as presented in .
The pressure has been calculated on the one-loop level for SU(2) and SU(3) Yang-Mills theory in ; we only consider the SU(2) case here. The pressure contains a temperature dependent ground-state contribution arising from caloron ’condensation’,
$$P_{\text{g.s.}}(\lambda )=2\lambda \mathrm{\Lambda }^4,$$
(4.2)
and a contribution associated with the one-loop diagrams shown in Fig. 4.1,
$$P_{\text{one-loop}}(\lambda )=\mathrm{\Lambda }^4\frac{2\lambda ^4}{(2\pi )^6}\left(\frac{2\pi ^4}{45}+6\overline{P}\left(4\pi e\lambda ^{3/2}\right)\right).$$
(4.3)
The fixed mass scale $`\mathrm{\Lambda }`$ is connected to the modulus of the composite Higgs field $`\left|\varphi \right|`$ by
$$\left|\varphi \right|^2=\frac{\mathrm{\Lambda }^3}{2\pi T},$$
(4.4)
the dimensionless temperature $`\lambda `$ is defined as
$$\lambda =\frac{2\pi T}{\mathrm{\Lambda }},$$
(4.5)
and the function $`\overline{P}`$ is given as
$$\overline{P}(a)=_0^{\mathrm{}}𝑑xx^2\mathrm{log}\left[1\mathrm{exp}\left(\sqrt{x^2+a^2}\right)\right].$$
(4.6)
There is also a ’nonthermal’ contribution $`\mathrm{\Delta }V`$. It is estimated as
$$\mathrm{\Delta }V<\frac{|\varphi |^4}{16\pi ^2}.$$
(4.7)
Because
$$\left|\frac{\mathrm{\Delta }V}{V}\right|<\frac{\lambda ^3}{32\pi ^2}<210^6,$$
(4.8)
it can be neglected in the electric phase. (Recall that $`\lambda >11.65`$ in the electric phase, cf. Sec. 3.1.)
#### 4.1.2 Feynman rules
For the calculation of the two-loop correction $`\mathrm{\Delta }P`$ to the pressure of SU(2) being in its electric phase, we have the equation (see )
$$\text{},$$
(4.9)
where double lines represent TLH modes and single lines stand for TLM modes. TLM modes will carry a color index 3, while the color indices 1 and 2 correspond to TLH modes. Eq. (4.9) is valid for SU(N). In the case $`\mathrm{N}=2`$ considered here, the second and third diagram do not occur.
The calculation of the thermodynamical pressure is performed in unitary gauge, where $`\varphi `$ is diagonal and the background is $`a_\mu ^{bg}=0`$. The remaining gauge freedom is used to gauge the TLM mode to transversality, $`_i\delta a_i^{TLM}=0`$ (Coulomb gauge). Unitary-Coulomb gauge is a completely fixed gauge, thus no Faddeev-Popov determinant needs to be considered and no ghost fields need to be introduced.
The computation is performed in the real-time formulation of thermal field theory. Spacetime is Minkowskian with signature $`(+,,,)`$. Matsubara sums, which originate when working in compactified imaginary time, are replaced with integrals over real Minkowskian momenta by means of contour integrals and analytic continuation. For our purposes, the real-time formalism is preferable because the implementation of constraints for the momenta of the participating fluctuations (Eqs. (4.15), (4.16)) is rather inconvenient in the imaginary-time formalism. Moreover, the quantum and thermal part of a given fluctuation can be clearly discerned in the real-time formalism.
In the real-time formulation and in unitary-Coulomb gauge, the Feynman rules employed in calculating the diagrams in Eq. (4.9) are as follows: The propagator for a free TLM-mode
$$D_{\mu \nu ,ab}^{TLM}(k,T)=\delta _{ab}P_{\mu \nu }^T(k)\left(\frac{i}{k^2+i\epsilon }+2\pi \delta (k^2)n_B(|k_0|/T)\right),$$
(4.10)
where<sup>1</sup><sup>1</sup>1 Static electric fields of long wavelength are completely screened due to the existence of an infinite real part in the Debye screening mass $`m_D=[\mathrm{\Pi }_{00}(k_0=0,𝒌0)]^{1/2}`$; cf. App. B
$$\begin{array}{cc}\hfill P_{00}^T(k)& =P_{0i}^T(k)=P_{i0}^T(k)=0,\hfill \\ \hfill P_{ij}^T(k)& =\delta _{ij}\frac{k_ik_j}{𝒌^2},\hfill \end{array}$$
(4.11)
and $`n_B`$ denotes the Bose-Einstein-distribution, $`n_B(x)=\frac{1}{e^x1}`$. The propagator for TLH-modes reads
$$D_{\mu \nu ,ab}^{TLH}(k,T)=\delta _{ab}\left(g_{\mu \nu }\frac{k_\mu k_\nu }{m^2}\right)\left(\frac{i}{k^2m^2+i\epsilon }+2\pi \delta (k^2m^2)n_B(|k_0|/T)\right).$$
(4.12)
The vertices are the usual ones (see ): The four-boson-vertex is
$$\begin{array}{c}\text{}=ie^2(2\pi )^4\delta ^4(p+q+r+s)[\epsilon _{fab}\epsilon _{fcd}(g^{\mu \rho }g^{\nu \sigma }g^{\mu \sigma }g^{\nu \rho })\hfill \\ \hfill +\epsilon _{fac}\epsilon _{fdb}(g^{\mu \sigma }g^{\rho \nu }g^{\mu \nu }g^{\rho \sigma })+\epsilon _{fad}\epsilon _{fbc}(g^{\mu \nu }g^{\sigma \rho }g^{\mu \rho }g^{\sigma \nu })],\end{array}$$
(4.13)
and the three-boson-vertex is
$$\text{}=e(2\pi )^4\delta ^4(p+q+k)\epsilon _{abc}[g^{\mu \nu }(qp)^\rho +g^{\nu \rho }(kq)^\mu +g^{\rho \mu }(pk)^\nu ].$$
(4.14)
According to Landsman and van Weert , one has to divide every diagram by $`i`$ and by the number of its vertices. Besides that, we demand that the momenta of quantum fluctuations are to be cut off at
$`\left|p^2m^2\right|\left|\varphi \right|^2`$ (4.15a)
in Minkowskian or
$`p_E^2+m^2\left|\varphi \right|^2`$ (4.15b)
in Euclidean signature because of the existence of the compositeness scale $`\left|\varphi \right|`$, cf. Sec. 3.1. The compositeness scale imposes a similar constraint on the center-of-mass energy flowing into the four-vertex,
$$\left|\left(p_1+p_2\right)^2\right|\left|\varphi \right|^2,$$
(4.16)
where $`p_i`$ are the ingoing momenta.
#### 4.1.3 Contributing diagrams
The various contributions to the pressure on the two-loop level will be denoted as follows: The first (last) diagram in Eq. (4.9) shall be referred to as $`\mathrm{\Delta }P^{\mathrm{MHH}}`$ ($`\mathrm{\Delta }P^{\mathrm{HH}}`$), since there are two TLH and one TLM particle (two TLH particles) involved. If we want to specify, for example, that in $`\mathrm{\Delta }P^{\mathrm{MH}}`$ the TLH fluctuation is a vacuum fluctuation and the TLM fluctuation a thermal fluctuation, we write $`\mathrm{\Delta }P_{\mathrm{TV}}^{\mathrm{MH}}`$ etc. The statistical factors of Eq. (4.9) are included in these definitions.
Which diagrams in Eq. (4.9) do actually contribute for SU(2)? Because of the structure of the vertex (4.14), the second and the third diagram in Eq. (4.9) vanish. Moreover, $`\mathrm{\Delta }P_{\mathrm{TTT}}^{\mathrm{MHH}}`$ vanishes (for any $`\mathrm{N}`$ and $`e`$), since the on-shell conditions $`p^2=k^2=m^2`$, $`q^2=0`$ and energy-momentum-conservation $`p+k+q=0`$ cannot be satisfied simultaneously.
As the mass of TLH fluctuations is connected to the gauge coupling via $`m^2=4e^2\left|\varphi \right|^2`$, the constraint (4.15) reads
$$\left|p^24e^2|\varphi |^2\right|\left|\varphi \right|^2\text{or}p_E^2\left(14e^2\right)\left|\varphi \right|^2.$$
(4.17)
This condition will eliminate all contributions containing quantum fluctuations of TLH modes, if the gauge coupling $`e`$ is larger than $`1/2`$. In this case, we only have the following nonvanishing contributions:
$$\begin{array}{cc}\hfill \mathrm{\Delta }P^{\mathrm{HH}}& =\mathrm{\Delta }P_{\mathrm{TT}}^{\mathrm{HH}}\hfill \\ \hfill \mathrm{\Delta }P^{\mathrm{MH}}& =\mathrm{\Delta }P_{\mathrm{TT}}^{\mathrm{MH}}+\mathrm{\Delta }P_{\mathrm{VT}}^{\mathrm{MH}}\hfill \\ \hfill \mathrm{\Delta }P^{\mathrm{MHH}}& =\mathrm{\Delta }P_{\mathrm{VTT}}^{\mathrm{MHH}}.\hfill \end{array}$$
(4.18)
So, using the above notation, Eq. (4.9) for $`e>1/2`$ simplifies to
$$\mathrm{\Delta }P=\mathrm{\Delta }P_{\mathrm{TT}}^{\mathrm{HH}}+\mathrm{\Delta }P_{\mathrm{TT}}^{\mathrm{MH}}+\mathrm{\Delta }P_{\mathrm{VT}}^{\mathrm{MH}}+\mathrm{\Delta }P_{\mathrm{VTT}}^{\mathrm{MHH}}.$$
(4.19)
### 4.2 Calculation
#### 4.2.1 Calculation of $`\mathrm{\Delta }P_{\mathrm{TT}}^{\mathrm{HH}}`$
We first write down the expression for $`\mathrm{\Delta }P_{\mathrm{TT}}^{\mathrm{HH}}`$ without imposing the kinematical constraints Eqs. (4.15), (4.16) and perform the constrained integrations afterwards still using the notation $`\mathrm{\Delta }P_{\mathrm{TT}}^{\mathrm{HH}}`$. According to the Feynman-rules given in section 4.1.2, the diagram shown in Fig. 4.2 with both fluctuations being thermal reads as:
$$\begin{array}{cc}\hfill \mathrm{\Delta }P_{\mathrm{TT}}^{\mathrm{HH}}& =\frac{1}{8}\frac{d^4k}{(2\pi )^4}\frac{d^4p}{(2\pi )^4}(e^2)[\epsilon _{adf}\epsilon _{fbc}(g^{\mu \nu }g^{\rho \sigma }g^{\mu \rho }g^{\nu \sigma })\hfill \\ & +\epsilon _{abf}\epsilon _{fdc}(g^{\mu \sigma }g^{\nu \rho }g^{\mu \rho }g^{\nu \sigma })+\epsilon _{acf}\epsilon _{fdb}(g^{\mu \sigma }g^{\nu \rho }g^{\mu \nu }g^{\sigma \rho })]\hfill \\ & \left(\delta _{ab}\right)\left(g_{\mu \nu }\frac{k_\mu k_\nu }{m^2}\right)2\pi \delta (k^2m^2)n_B(|k_0|/T)\hfill \\ & \left(\delta _{cd}\right)\left(g_{\rho \sigma }\frac{p_\rho p_\sigma }{m^2}\right)2\pi \delta (p^2m^2)n_B(|p_0|/T)\hfill \\ & =\frac{e^2}{2}\frac{d^4k}{(2\pi )^4}\frac{d^4p}{(2\pi )^4}\left(123\frac{p^2}{m^2}3\frac{k^2}{m^2}+\frac{p^2k^2}{m^4}\frac{(pk)^2}{m^4}\right)\hfill \\ & 2\pi \delta (k^2m^2)n_B(|k_0|/T)2\pi \delta (p^2m^2)n_B(|p_0|/T).\hfill \end{array}$$
(4.20)
(The contraction of Lorentz and color indices is deferred to Appendix B.1.) In evaluating this integral, we will first perform the integrations over the zero-components of the momenta and thus eliminate the $`\delta `$-functions. Therefore we rewrite the $`\delta `$-functions as follows:
$$\begin{array}{cc}\hfill \mathrm{\Delta }P_{\mathrm{TT}}^{\mathrm{HH}}& =\frac{e^2}{2}\frac{d^3k}{(2\pi )^3}\frac{d^3p}{(2\pi )^3}𝑑k_0𝑑p_0\left(7\frac{p_0^2k_0^22p_0k_0𝒑𝒌+(𝒑𝒌)^2}{m^4}\right)\hfill \\ & \frac{1}{2\sqrt{𝒌^2+m^2}}\left[\delta \left(k_0\sqrt{𝒌^2+m^2}\right)+\delta \left(k_0+\sqrt{𝒌^2+m^2}\right)\right]n_B(|k_0|/T)\hfill \\ & \frac{1}{2\sqrt{𝒑^2+m^2}}\left[\delta \left(p_0\sqrt{𝒑^2+m^2}\right)+\delta \left(p_0+\sqrt{𝒑^2+m^2}\right)\right]n_B(|p_0|/T).\hfill \end{array}$$
(4.21)
Now we notice that the two terms containing
$$\delta \left(k_0\sqrt{𝒌^2+m^2}\right)\delta \left(p_0\sqrt{𝒑^2+m^2}\right)$$
(4.22)
or
$$\delta \left(k_0+\sqrt{𝒌^2+m^2}\right)\delta \left(p_0+\sqrt{𝒑^2+m^2}\right),$$
(4.23)
respectively, yield the same result since the rest of the integrand is invariant under simultaneous reflections $`p_0p_0`$, $`k_0k_0`$. The same holds true for the remaining two summands. So the expression is split into two contributions as
$$\begin{array}{cc}\hfill \mathrm{\Delta }P_{\mathrm{TT}}^{\mathrm{HH}}& =\frac{e^2}{4}\underset{\pm }{}\frac{d^3k}{(2\pi )^3}\frac{d^3p}{(2\pi )^3}\frac{1}{\sqrt{𝒌^2+m^2}\sqrt{𝒑^2+m^2}}\hfill \\ & \left(7\frac{(𝒑^2+m^2)(𝒌^2+m^2)2\sqrt{𝒑^2+m^2}\sqrt{𝒌^2+m^2}𝒑𝒌+(𝒑𝒌)^2}{m^4}\right)\hfill \\ & n_B\left(\sqrt{𝒑^2+m^2}/T\right)n_B\left(\sqrt{𝒌^2+m^2}/T\right)\hfill \\ & =\frac{e^2}{4}\underset{\pm }{}\frac{d^3k}{(2\pi )^3}\frac{d^3p}{(2\pi )^3}\frac{1}{\sqrt{𝒌^2+m^2}\sqrt{𝒑^2+m^2}}\hfill \\ & \left(6\frac{𝒑^2}{m^2}\frac{𝒌^2}{m^2}\frac{𝒑^2𝒌^2+(𝒑𝒌)^2}{m^4}\pm 2𝒑𝒌\frac{\sqrt{𝒑^2+m^2}\sqrt{𝒌^2+m^2}}{m^4}\right)\hfill \\ & n_B\left(\sqrt{𝒑^2+m^2}/T\right)n_B\left(\sqrt{𝒌^2+m^2}/T\right).\hfill \end{array}$$
(4.24)
The tree-level mass of TLH-modes in SU(2) is given as
$$m=2e|\varphi |,$$
(4.25)
so the introduction of dimensionless momenta $`𝒚`$ and $`𝒙`$ according to
$$𝒌=𝒚|\varphi |\text{and}𝒑=𝒙|\varphi |$$
(4.26)
is suggested. Moreover, temperature $`T`$ and compositeness scale $`\left|\varphi \right|`$ can be expressed in terms of the dimensionless temperature $`\lambda `$ and the mass scale $`\mathrm{\Lambda }`$ (see Eqs. (4.4) and (4.5)) as
$$\begin{array}{cc}\hfill |\varphi |& =\mathrm{\Lambda }\lambda ^{1/2},\hfill \\ \hfill |\varphi |/T& =2\pi \lambda ^{3/2}.\hfill \end{array}$$
(4.27)
Finally, we introduce three-dimensional polar coordinates for the scaled momenta. The angle between $`𝒙`$ and $`𝒚`$ can be chosen to be the polar angle $`\theta `$. $`\mathrm{\Delta }P_{\mathrm{TT}}^{\mathrm{HH}}`$ is then expressed as
$$\begin{array}{cc}\hfill \mathrm{\Delta }P_{\mathrm{TT}}^{\mathrm{HH}}& =\frac{e^2\mathrm{\Lambda }^4\lambda ^2}{2(2\pi )^4}\underset{\pm }{}_0^{\mathrm{}}𝑑x_0^{\mathrm{}}𝑑y_1^1d\mathrm{cos}\theta \frac{x^2y^2}{\sqrt{x^2+4e^2}\sqrt{y^2+4e^2}}\hfill \\ & \left(6\frac{x^2}{4e^2}\frac{y^2}{4e^2}\frac{x^2y^2}{16e^4}\left(1+\mathrm{cos}^2\theta \right)\pm 2xy\mathrm{cos}\theta \frac{\sqrt{x^2+4e^2}\sqrt{y^2+4^2}}{16e^4}\right)\hfill \\ & n_B\left(2\pi \lambda ^{3/2}\sqrt{x^2+4e^2}\right)n_B\left(2\pi \lambda ^{3/2}\sqrt{y^2+4e^2}\right).\hfill \end{array}$$
(4.28)
Now we have to implement the constraint for the center-of-mass energy in the vertex,
$$\left|(p+k)^2\right|\left|\varphi \right|^2.$$
(4.29)
Of course, this condition has to undergo the same manipulations as the integrand itself: $`p^2`$ and $`k^2`$ have to be replaced with $`m^2`$, $`p_0k_0`$ with $`\pm \sqrt{𝒑^2+m^2}\sqrt{𝒌^2+m^2}`$,
$$\left|2m^2\pm 2\sqrt{𝒑^2+m^2}\sqrt{𝒌^2+m^2}2𝒑𝒌\right||\varphi |^2,$$
(4.30)
the momenta have to be rescaled as Eq. (4.26), the expression for the mass from Eq. (4.25) has to be inserted, and a change to polar coordinates has to be done. We then have
$$\left|4e^2\pm \sqrt{x^2+4e^2}\sqrt{y^2+4e^2}xy\mathrm{cos}\theta \right|\frac{1}{2}.$$
(4.31)
For $`e`$ larger than $`\frac{1}{2\sqrt{2}}`$, condition (4.31) with the positive sign cannot be fulfilled, and the corresponding integral in Eq. (4.28) is zero. For the negative sign, condition (4.31) can be fulfilled, and we will implement it as an additional constraint on the region of integration for $`\mathrm{cos}\theta `$. The condition
$$\left|4e^2\sqrt{x^2+4e^2}\sqrt{y^2+4e^2}xy\mathrm{cos}\theta \right|=4e^2+\sqrt{x^2+4e^2}\sqrt{y^2+4e^2}+xy\mathrm{cos}\theta \frac{1}{2}$$
(4.32)
is equivalent to
$$\mathrm{cos}\theta \frac{1}{xy}\left(\frac{1}{2}+4e^2\sqrt{x^2+4e^2}\sqrt{y^2+4e^2}\right)g(x,y),$$
(4.33)
and thus $`\mathrm{cos}\theta `$ is to be integrated in the range $`[1,1](\mathrm{},g]`$. This can be done by replacing the upper limit of integration for $`\mathrm{cos}\theta `$ with
$$\mathrm{max}\{1,\mathrm{min}\{1,g\}\}.$$
(4.34)
So, we have in total
$$\begin{array}{cc}\hfill \mathrm{\Delta }P_{\mathrm{TT}}^{\mathrm{HH}}& =\frac{e^2\mathrm{\Lambda }^4\lambda ^2}{2(2\pi )^4}_0^{\mathrm{}}𝑑x_0^{\mathrm{}}𝑑y\underset{1}{\overset{\mathrm{max}\{1,\mathrm{min}\{1,g\}\}}{}}d\mathrm{cos}\theta \frac{x^2y^2}{\sqrt{x^2+4e^2}\sqrt{y^2+4e^2}}\hfill \\ & \left(6\frac{x^2}{4e^2}\frac{y^2}{4e^2}\frac{x^2y^2}{16e^4}\left(1+\mathrm{cos}^2\theta \right)2xy\mathrm{cos}\theta \frac{\sqrt{x^2+4e^2}\sqrt{y^2+4^2}}{16e^4}\right)\hfill \\ & n_B\left(2\pi \lambda ^{3/2}\sqrt{x^2+4e^2}\right)n_B\left(2\pi \lambda ^{3/2}\sqrt{y^2+4e^2}\right).\hfill \end{array}$$
(4.35)
In this form the integrals can be performed numerically. Instead of doing this, we will have a closer look at the function $`g`$ in order to be able to give the boundaries of integration explicitly. For the following compare with Fig. 4.3.
If $`g`$ is greater than 1, Eq. (4.33) is fulfilled automatically and $`\mathrm{cos}\theta `$ has to be integrated from -1 to 1; this is the case in the white area in Fig. 4.3, bounded by the curve
$$b(x)=\frac{x8e^2x+\sqrt{1+16e^2}\sqrt{x^2+4e^2}}{8e^2}.$$
(4.36)
This curve intersects with the $`x`$-axis at
$$a=\sqrt{1+\frac{1}{16e^2}}.$$
(4.37)
If $`1<g<1`$, the upper limit for the integration of $`\mathrm{cos}\theta `$ is $`g`$. The region where this is true is shaded in grey in Fig. 4.3; it is enclosed by the lines
$$\begin{array}{cc}\hfill c(x)& =\frac{x+8e^2x\sqrt{1+16e^2}\sqrt{x^2+4e^2}}{8e^2},\hfill \\ \hfill d(x)& =\frac{x+8e^2x+\sqrt{1+16e^2}\sqrt{x^2+4e^2}}{8e^2},\hfill \end{array}$$
(4.38)
and $`b(x)`$ as given above. For all other values of $`x`$ and $`y`$ (shaded in black), we have $`g<1`$ and thus no contribution. This can be summarized as follows: In Eq. (4.35) the integral $`𝑑x𝑑yd\mathrm{cos}\theta `$ is decomposed as
$$\begin{array}{c}\hfill _0^a𝑑x_0^b𝑑y_1^1d\mathrm{cos}\theta +_0^a𝑑x_b^d𝑑y_1^gd\mathrm{cos}\theta +_a^{\mathrm{}}𝑑x_c^d𝑑y_1^gd\mathrm{cos}\theta \end{array}$$
(4.39)
where integrand and prefactor are as in Eq. (4.35). The integrations over $`\mathrm{cos}\theta `$ can be performed analytically; further evaluation needs to be done numerically. For the result see Fig. 4.6.
Fig. 4.3 illustrates that the integration over loop momenta is strongly restricted to a narrow band around $`|𝒑|=|𝒌|`$ due to the compositeness constraint Eq. (4.16).
#### 4.2.2 Calculation of $`\mathrm{\Delta }P_{\mathrm{TT}}^{\mathrm{MH}}`$ and $`\mathrm{\Delta }P_{\mathrm{VT}}^{\mathrm{MH}}`$
For the diagram depicted in Fig. 4.4, the Feynman rules yield
$`\mathrm{\Delta }P^{\mathrm{MH}}`$ $`={\displaystyle \frac{1}{8}}{\displaystyle }{\displaystyle \frac{d^4k}{(2\pi )^4}}{\displaystyle }{\displaystyle \frac{d^4p}{(2\pi )^4}}(e^2)[\epsilon _{adf}\epsilon _{fbc}(g^{\mu \nu }g^{\rho \sigma }g^{\mu \rho }g^{\nu \sigma })`$
$`+\epsilon _{abf}\epsilon _{fdc}(g^{\mu \sigma }g^{\nu \rho }g^{\mu \rho }g^{\nu \sigma })+\epsilon _{acf}\epsilon _{fdb}(g^{\mu \sigma }g^{\nu \rho }g^{\mu \nu }g^{\sigma \rho })]`$
$`\left(\delta _{cd}\right)\left(g_{\rho \sigma }{\displaystyle \frac{p_\rho p_\sigma }{m^2}}\right)\left[{\displaystyle \frac{i}{p^2m^2+i\epsilon }}+2\pi \delta (p^2m^2)n_B(|p_0|/T)\right]`$ (4.40)
$`\left(\delta _{ab}\right)P_{\mu \nu }^T(k)\left[{\displaystyle \frac{i}{k^2+i\epsilon }}+2\pi \delta (k^2)n_B(|k_0|/T)\right]`$
$`={\displaystyle \frac{e^2}{2}}{\displaystyle \frac{d^4k}{(2\pi )^4}\frac{d^4p}{(2\pi )^4}\left(6+2\frac{p^2}{m^2}+\frac{𝒑^2}{m^2}\frac{(𝒑𝒌)^2}{m^2𝒌^2}\right)}`$
$`\left[{\displaystyle \frac{i}{p^2m^2+i\epsilon }}+2\pi \delta (p^2m^2)n_B(|p_0|/T)\right]\left[{\displaystyle \frac{i}{k^2+i\epsilon }}+2\pi \delta (k^2)n_B(|k_0|/T)\right].`$
##### Calculation of $`\mathrm{\Delta }P_{\mathrm{TT}}^{\mathrm{MH}}`$
For both fluctuations being thermal, we have
$$\begin{array}{cc}\hfill \mathrm{\Delta }P_{\mathrm{TT}}^{\mathrm{MH}}& =\frac{e^2}{2}\frac{d^4k}{(2\pi )^4}\frac{d^4p}{(2\pi )^4}\left(6+2\frac{p^2}{m^2}+\frac{𝒑^2}{m^2}\frac{(𝒑𝒌)^2}{m^2𝒌^2}\right)\hfill \\ & 2\pi \delta (k^2)n_B(|k_0|/T)2\pi \delta (p^2m^2)n_B(|p_0|/T).\hfill \end{array}$$
(4.41)
The evaluation of this diagram is similar to the calculation in Sec. 4.2.1. Integrating over the zero components of $`p`$ and $`k`$, we have
$$\begin{array}{cc}\hfill \mathrm{\Delta }P_{\mathrm{TT}}^{\mathrm{MH}}& =\frac{e^2}{2}\frac{d^3k}{(2\pi )^3}\frac{d^3p}{(2\pi )^3}𝑑k_0𝑑p_0\left(4+2\frac{𝒑^2}{m^2}\frac{(𝒑𝒌)^2}{m^2𝒌^2}\right)\hfill \\ & \frac{1}{2\sqrt{𝒑^2+m^2}}\left[\delta \left(p_0\sqrt{𝒑^2+m^2}\right)+\delta \left(p_0+\sqrt{𝒑^2+m^2}\right)\right]n_B(|p_0|/T)\hfill \\ & \frac{1}{2\sqrt{𝒌^2}}\left[\delta \left(k_0\sqrt{𝒌^2}\right)+\delta \left(k_0+\sqrt{𝒌^2}\right)\right]n_B(|k_0|/T)\hfill \\ & =\frac{e^2}{2}\underset{\pm }{}\frac{d^3k}{(2\pi )^3}\frac{d^3p}{(2\pi )^3}\left(4+2\frac{𝒑^2}{m^2}\frac{(𝒑𝒌)^2}{m^2𝒌^2}\right)\hfill \\ & \frac{1}{2\sqrt{𝒑^2+m^2}\sqrt{𝒌^2}}n_B\left(\sqrt{𝒌^2}/T\right)n_B\left(\sqrt{𝒑^2+m^2}/T\right).\hfill \end{array}$$
(4.42)
Inserting the mass given in Eq. (4.25), scaling momenta as in Eq. (4.26), introducing dimensionless variables as in Eq. (4.27) and polar coordinates, we obtain
$$\begin{array}{cc}\hfill \mathrm{\Delta }P_{\mathrm{TT}}^{\mathrm{MH}}& =\frac{e^2\mathrm{\Lambda }^4\lambda ^2}{2(2\pi )^4}\underset{\pm }{}_0^{\mathrm{}}𝑑x_0^{\mathrm{}}𝑑y_1^1d\mathrm{cos}\theta \hfill \\ & \left(4+\frac{x^2}{4e^2}\frac{x^2\mathrm{cos}^2\theta }{4e^2}\right)\frac{x^2y}{\sqrt{x^2+4e^2}}\hfill \\ & n_B\left(2\pi \lambda ^{3/2}\sqrt{x^2+4e^2}\right)n_B\left(2\pi \lambda ^{3/2}y\right).\hfill \end{array}$$
(4.43)
At this stage the $`\pm `$ contributions in this equation are the same, but after imposing the condition
$$\left|(p+k)^2\right|\left|\varphi \right|^2,$$
(4.44)
this is no longer the case. By the above manipulations, Eq. (4.44) is transformed into
$$\left|m^2\pm 2\sqrt{𝒑^2+m^2}\sqrt{𝒌^2}2𝒑𝒌\right|\left|\varphi \right|^2$$
(4.45)
and finally into
$$\left|4e^2\pm 2y\sqrt{x^2+4e^2}2xy\mathrm{cos}\theta \right|1.$$
(4.46)
Again, condition (4.46), which limits the range of integration for $`\mathrm{cos}\theta `$, can only be fulfilled for the negative sign. The condition
$$\left|4e^22y\sqrt{x^2+4e^2}2xy\mathrm{cos}\theta \right|1$$
(4.47)
is satisfied if and only if
$$\frac{1+4e^22y\sqrt{x^2+4e^2}}{2xy}\tau _2\mathrm{cos}\theta \tau _1\frac{+1+4e^22y\sqrt{x^2+4e^2}}{2xy},$$
(4.48)
and thus $`\mathrm{cos}\theta `$ is to be integrated over the interval $`[1,1][\tau _2,\tau _1]`$. So the upper and lower limits of integration in Eq. (4.43) are
$$\mathrm{max}\{1,\mathrm{min}\{1,\tau _1\}\}\text{and}\mathrm{min}\{1,\mathrm{max}\{1,\tau _2\}\}$$
(4.49)
respectively. Hence the result is
$$\begin{array}{cc}\hfill \mathrm{\Delta }P_{\mathrm{TT}}^{\mathrm{MH}}& =\frac{e^2\mathrm{\Lambda }^4\lambda ^2}{2(2\pi )^4}_0^{\mathrm{}}𝑑x_0^{\mathrm{}}𝑑y\underset{\mathrm{min}\{1,\mathrm{max}\{1,\tau _2\}\}}{\overset{\mathrm{max}\{1,\mathrm{min}\{1,\tau _1\}\}}{}}d\mathrm{cos}\theta \frac{x^2y}{\sqrt{x^2+4e^2}}\hfill \\ & \left(4+\frac{x^2}{4e^2}\frac{x^2\mathrm{cos}^2\theta }{4e^2}\right)n_B\left(2\pi \lambda ^{3/2}\sqrt{x^2+4e^2}\right)n_B\left(2\pi \lambda ^{3/2}y\right).\hfill \end{array}$$
(4.50)
Again, we want to give the limits of integration explicitly. To do this, the points with $`\tau _{1/2}=\pm 1`$ need to be determined. They are
$$\begin{array}{cc}\hfill \tau _1=\pm 1& y=y_{1\pm }\frac{4e^2+1}{2\left(\pm x+\sqrt{x^2+4e^2}\right)}\hfill \\ \hfill \tau _2=\pm 1& y=y_{2\pm }\frac{4e^21}{2\left(\pm x+\sqrt{x^2+4e^2}\right)}.\hfill \end{array}$$
(4.51)
The following inequalities are satisfied:
$$\begin{array}{cc}\hfill y_{2+}y_2y_{1+}y_1& x\eta \frac{2e}{\sqrt{16e^21}}\hfill \\ \hfill y_{2+}y_{1+}y_2y_1& x\eta .\hfill \end{array}$$
(4.52)
Taking into account all possible combinations of $`x`$ and $`y`$ and the corresponding limits of integration, the integral in Eq. (4.50) splits into six parts, namely
$$\begin{array}{cc}& _0^\eta 𝑑x_{y_{2+}}^{y_2}𝑑y_{\tau _2}^1𝑑t+_0^\eta 𝑑x_{y_2}^{y_{1+}}𝑑y_1^1𝑑t+_0^\eta 𝑑x_{y_{1+}}^{y_1}𝑑y_1^{\tau _1}𝑑t\hfill \\ & +_\eta ^{\mathrm{}}𝑑x_{y_{2+}}^{y_{1+}}𝑑y_{\tau _2}^1𝑑t+_\eta ^{\mathrm{}}𝑑x_{y_{1+}}^{y_2}𝑑y_{\tau _2}^{\tau _1}𝑑t+_\eta ^{\mathrm{}}𝑑x_{y_2}^{y_1}𝑑y_1^{\tau _1}𝑑t.\hfill \end{array}$$
(4.53)
The result of the numerical evaluation is shown in Fig. 4.7.
##### Calculation of $`\mathrm{\Delta }P_{\mathrm{VT}}^{\mathrm{MH}}`$
For this case, we have from Eq. (4.2.2) the expression
$$\begin{array}{cc}\hfill \mathrm{\Delta }P_{\mathrm{VT}}^{\mathrm{MH}}& =\frac{e^2}{2}\frac{d^4k}{(2\pi )^4}\frac{d^4p}{(2\pi )^4}\mathrm{\hspace{0.17em}2}\pi \delta \left(p^2m^2\right)n_B\left(\left|p_0\right|/T\right)\hfill \\ & \left(6+2\frac{p^2}{m^2}+\frac{𝒑^2}{m^2}\frac{(𝒑𝒌)^2}{m^2𝒌^2}\right)\frac{i}{k^2+i\epsilon }.\hfill \end{array}$$
(4.54)
As there is one thermal fluctuation ($`p`$) and one vacuum fluctuation ($`k`$) involved, only $`p_0`$ can be integrated by eliminating the $`\delta `$-function,
$$\begin{array}{cc}\hfill \mathrm{\Delta }P_{\mathrm{VT}}^{\mathrm{MH}}& =\frac{e^2}{2}\frac{d^4k}{(2\pi )^4}\frac{d^3p}{(2\pi )^3}𝑑p_0n_B\left(\sqrt{𝒑^2+m^2}/T\right)\frac{i}{k^2+i\epsilon }\hfill \\ & \left(4+\frac{𝒑^2}{m^2}\frac{(𝒑𝒌)^2}{m^2𝒌^2}\right)\frac{\delta \left(p_0\sqrt{𝒑^2+m^2}\right)+\delta \left(p_0+\sqrt{𝒑^2+m^2}\right)}{2\sqrt{𝒑^2+m^2}}\hfill \\ & =\frac{e^2}{2}\underset{\pm }{}\frac{d^4k}{(2\pi )^4}\frac{d^3p}{(2\pi )^3}n_B\left(\sqrt{𝒑^2+m^2}/T\right)\frac{i}{k^2+i\epsilon }\hfill \\ & \left(4+\frac{𝒑^2}{m^2}\frac{(𝒑𝒌)^2}{m^2𝒌^2}\right)\frac{1}{2\sqrt{𝒑^2+m^2}}.\hfill \end{array}$$
(4.55)
The implementation of the compositeness constraints
$$\left|(p+k)^2\right|\left|\varphi \right|^2$$
(4.56)
and
$$\left|k^2\right|\left|\varphi \right|^2$$
(4.57)
is more difficult as in the previous calculations. We will therefore ignore Eq. (4.56) and give an estimate for $`\mathrm{\Delta }P_{\mathrm{VT}}^{\mathrm{MH}}`$ by only taking into account (4.57). The two contributions $`\pm `$ are equal, $`k`$ is analytically continued to Euclidean momenta. This yields the upper bound
$$\begin{array}{cc}\hfill \left|\mathrm{\Delta }P_{\mathrm{VT}}^{\mathrm{MH}}\right|& \frac{e^2}{2}\underset{k_E^2\left|\varphi \right|^2}{}\frac{d^4k_E}{(2\pi )^4}\frac{d^3p}{(2\pi )^3}\frac{n_B\left(\sqrt{𝒑^2+m^2}/T\right)}{\sqrt{𝒑^2+m^2}}\frac{1}{k_E^2}\left|4+\frac{𝒑^2}{m^2}\frac{(𝒑𝒌)^2}{m^2𝒌^2}\right|.\hfill \end{array}$$
(4.58)
As usual (inserting the expression for the mass, scaling the momenta and polar coordinates for the 3-vector $`𝒑`$ and the 4-vector $`k_E`$), we obtain
$$\begin{array}{cc}\hfill \left|\mathrm{\Delta }P_{\mathrm{VT}}^{\mathrm{MH}}\right|& \frac{e^2\mathrm{\Lambda }^4\lambda ^2}{8(2\pi )^4}_0^{\mathrm{}}𝑑x_1^1d\mathrm{cos}\theta \frac{x^2}{\sqrt{x^2+4e^2}}\hfill \\ & n_B\left(2\pi \lambda ^{3/2}\sqrt{x^2+4e^2}\right)\left|4+\frac{x^2}{4e^2}(1\mathrm{cos}^2\theta )\right|.\hfill \end{array}$$
(4.59)
The numerical evaluation of this upper bound is shown in Fig. 4.8.
#### 4.2.3 Calculation of $`\mathrm{\Delta }P_{\mathrm{VTT}}^{\mathrm{MHH}}`$
The contribution $`\mathrm{\Delta }P_{\mathrm{VTT}}^{\mathrm{MHH}}`$ shown in Fig. 4.5 corresponds to the expression
$$\begin{array}{cc}\hfill \mathrm{\Delta }P_{\mathrm{VTT}}^{\mathrm{MHH}}& =\frac{1}{4}\frac{e^2}{2i}\frac{d^4k}{(2\pi )^4}\frac{d^4p}{(2\pi )^4}\frac{d^4q}{(2\pi )^4}(2\pi )^4\delta (p+q+k)\hfill \\ & \epsilon _{acf}\left[g^{\mu \rho }(pk)^\lambda +g^{\rho \lambda }(kq)^\mu +g^{\lambda \mu }(qp)^\rho \right]\hfill \\ & \epsilon _{bdg}\left[g^{\nu \sigma }(kp)^\kappa +g^{\sigma \kappa }(qk)^\nu +g^{\kappa \nu }(pq)^\sigma \right]\hfill \\ & (\delta _{ab})\left(g_{\mu \nu }\frac{p_\mu p_\nu }{m^2}\right)2\pi \delta (p^2m^2)n_B(|p_0|/T)\hfill \\ & (\delta _{cd})\left(g_{\rho \sigma }\frac{k_\rho k_\sigma }{m^2}\right)2\pi \delta (k^2m^2)n_B(|k_0|/T)(\delta _{fg})P_{\lambda \kappa }^T(q)\frac{i}{q^2+i\epsilon }\hfill \\ & =\frac{e^2}{4}\frac{d^4p}{(2\pi )^4}\frac{d^4k}{(2\pi )^4}\frac{1}{(k+p)^2+i\epsilon }\hfill \\ & \left[16\left(m^2\frac{(kp)^2}{m^2}\right)\frac{𝒌^2𝒑^2(𝒌𝒑)^2}{(𝒑+𝒌)^2}\left(8+4\frac{(kp)^2}{m^4}\right)\right]\hfill \\ & 2\pi \delta \left(p^2m^2\right)n_B\left(\left|p_0\right|/T\right)2\pi \delta \left(k^2m^2\right)n_B\left(\left|k_0\right|/T\right).\hfill \end{array}$$
(4.60)
The vacuum propagator $`[(k+p)^2+i\epsilon ]^1`$ has a singularity at $`(p+k)^2=2(m^2+kp)=0`$, that is at $`kp=m^2`$. In spite of this, the term containing
$$16\left(m^2\frac{(kp)^2}{m^2}\right)$$
(4.61)
is regular everywhere even for $`\epsilon =0`$ because
$$\frac{16}{(k+p)^2}\frac{m^4(kp)^2}{m^2}=\frac{16}{m^2}\frac{m^4(kp)^2}{2(m^2+kp)}=\frac{8}{m^2}(m^2kp)=8\left(1\frac{kp}{m^2}\right).$$
(4.62)
So $`i\epsilon `$ can be dropped for this term. This is, however, not the case for the second part which is proportional to
$$\frac{𝒌^2𝒑^2(𝒌𝒑)^2}{(𝒑+𝒌)^2}\left(8+4\frac{(kp)^2}{m^4}\right).$$
(4.63)
In this diagram we again have to take care of constraining the momenta of vacuum fluctuations,
$$\left|q^2\right|=\left|(p+k)^2\right|=2\left|m^2+pk\right|\left|\varphi \right|^2.$$
(4.64)
This is exactly the same condition as has been found in the calculation of $`\mathrm{\Delta }P_{\mathrm{TT}}^{\mathrm{HH}}`$ (compare Eq. (4.29) and the following), so we can take the integration limits from Sec. 4.2.1.
We once again reformulate the $`\delta `$-functions,
$$\begin{array}{cc}\hfill \mathrm{\Delta }P_{\mathrm{VTT}}^{\mathrm{MHH}}& =\frac{e^2}{4}\frac{d^3p}{(2\pi )^3}\frac{d^3k}{(2\pi )^3}𝑑p_0𝑑k_0\frac{1}{2(m^2+kp)+i\epsilon }\hfill \\ & \frac{1}{2\sqrt{𝒑^2+m^2}}\left[\delta \left(p_0\sqrt{𝒑^2+m^2}\right)+\delta \left(p_0+\sqrt{𝒑^2+m^2}\right)\right]n_B\left(\left|p_0\right|/T\right)\hfill \\ & \frac{1}{2\sqrt{𝒌^2+m^2}}\left[\delta \left(k_0\sqrt{𝒌^2+m^2}\right)+\delta \left(k_0+\sqrt{𝒌^2+m^2}\right)\right]n_B\left(\left|k_0\right|/T\right)\hfill \\ & \left[16\left(m^2\frac{(kp)^2}{m^2}\right)\frac{𝒌^2𝒑^2(𝒌𝒑)^2}{(𝒑+𝒌)^2}\left(8+4\frac{(kp)^2}{m^4}\right)\right]\hfill \end{array}$$
(4.65)
and perform the integrations over $`p_0`$ and $`k_0`$,
$$\begin{array}{cc}\hfill \mathrm{\Delta }P_{\mathrm{VTT}}^{\mathrm{MHH}}& =\frac{e^2}{16}\frac{d^3p}{(2\pi )^3}\frac{d^3k}{(2\pi )^3}\frac{1}{m^2\sqrt{𝒑^2+m^2}\sqrt{𝒌^2+m^2}+𝒑𝒌+i\epsilon }\hfill \\ & \frac{1}{\sqrt{𝒌^2+m^2}\sqrt{𝒑^2+m^2}}n_B\left(\sqrt{𝒑^2+m^2}/T\right)n_B\left(\sqrt{𝒌^2+m^2}/T\right)\hfill \\ & [16(m^2\frac{\left(\sqrt{𝒌^2+m^2}\sqrt{𝒑^2+m^2}+𝒌𝒑\right)^2}{m^2})\hfill \\ & \frac{𝒌^2𝒑^2(𝒌𝒑)^2}{(𝒑+𝒌)^2}(8+4\frac{\left(\sqrt{𝒌^2+m^2}\sqrt{𝒑^2+m^2}+𝒌𝒑\right)^2}{m^4})],\hfill \end{array}$$
(4.66)
rescale momenta by setting
$$𝒑=𝒙\left|\varphi \right|\text{and}𝒌=𝒚\left|\varphi \right|,$$
(4.67)
go to polar coordinates, and introduce dimensionless variables $`\lambda `$ and $`\mathrm{\Lambda }`$ as in Eq. (4.27). Finally we arrive at
$$\begin{array}{cc}\hfill \mathrm{\Delta }P_{\mathrm{VTT}}^{\mathrm{MHH}}& =\frac{e^2\mathrm{\Lambda }^4\lambda ^2}{8(2\pi )^4}𝑑x𝑑y𝑑t\frac{x^2y^2}{\sqrt{x^2+4e^2}\sqrt{y^2+4e^2}}\hfill \\ & n_B\left(2\pi \lambda ^{3/2}\sqrt{x^2+4e^2}\right)n_B\left(2\pi \lambda ^{3/2}\sqrt{y^2+4e^2}\right)\hfill \\ & \frac{1}{4e^2\sqrt{x^2+4e^2}\sqrt{y^2+4e^2}xyt+i\epsilon }[\frac{x^2y^2(t^21)}{x^2+y^2+2xyt}\hfill \\ & \left(12+\frac{x^2}{e^2}+\frac{y^2}{e^2}+\frac{x^2y^2}{4e^4}(1+t^2)+\frac{xyt}{2e^4}\sqrt{x^2+4e^2}\sqrt{y^2+4e^2}\right)\hfill \\ & 16(x^2+y^2+\frac{x^2y^2}{4e^2}(1+t^2)+\frac{xyt}{2e^2}\sqrt{x^2+4e^2}\sqrt{y^2+4e^2})].\hfill \end{array}$$
(4.68)
The condition in Eq. (4.64) is implemented exactly as in the case of $`\mathrm{\Delta }P_{\mathrm{TT}}^{\mathrm{HH}}`$, see Eq. (4.36) to (4.39); $`\mathrm{cos}\theta `$ has been abbreviated by $`t`$ in Eq. (4.68).
As already indicated above, the term proportional to $`\frac{x^2y^2(t^21)}{x^2+y^2+2xyt}`$ has a singularity. We evaluate the integral numerically by prescribing a small value for $`\epsilon `$. In Tab. 4.1, one can see that for decreasing $`\epsilon `$ the real part converges to a finite value, while the imaginary part converges to zero. For our purposes, working with $`\epsilon =10^8`$ is sufficient for determining the real part; the imaginary part is ignored. The second term is regular, and the evaluation is straightforward by setting $`\epsilon =0`$. The result of numerical evaluation is shown in Fig. 4.9.
### 4.3 Results
To compare the corrections arising from the two-loop diagrams to the one-loop pressure, we plot the ratio of each of the contributions to $`\mathrm{\Delta }P`$ and $`P_{\text{one-loop}}`$ from Eq. (4.3) as a function of the dimensionless temperature $`\lambda _c=11.65\lambda 250`$. The results are shown in Fig. 4.6 through Fig. 4.9. The one-loop pressure does not include the ground-state contribution Eq. (4.2). For the effective gauge coupling, the plateau value $`e=5.1`$ is used for all temperatures $`\lambda `$. Throughout most of the electric phase, this is admissible, but the logarithmic pole of $`e`$ at the critical temperature $`\lambda _c`$ is ignored. The nonlocal diagram is the dominating contribution for $`\lambda <100`$. Throughout the electric phase, the corrections arising from two-loop diagrams are tiny; the ratio of two-loop to one-loop contribution is $`210^3`$ at most. With rising temperature, the two-loop contributions decrease. The contributions $`\mathrm{\Delta }P_{\mathrm{VTT}}^{\mathrm{MHH}}`$ and $`\mathrm{\Delta }P_{\mathrm{TT}}^{\mathrm{HH}}`$ approach zero for large temperatures. In contrast to that, $`\mathrm{\Delta }P_{\mathrm{TT}}^{\mathrm{MH}}`$ becomes constant. For $`\mathrm{\Delta }P_{\mathrm{VT}}^{\mathrm{MH}}`$ we obtained only a rough estimate. For temperatures close to the phase transition, there is a dip in the two-loop corrections to the pressure. The minimum of the dominating contribution is at $`\lambda 27`$.
The microscopic interpretation is as follows: Close to the phase transition at $`\lambda _c`$, the monopole mass $`M_{\text{monopole}}\frac{T}{e}`$ decreases sizeably. This increases the scattering of TLM modes off magnetic monopoles, with the consequence of a dip in the dominating two-loop contribution. With rising temperature the monopoles become massive and dilute, and scattering processes are suppressed. But as even for asymptotically high temperatures massive and dilute scattering centers are present, the contribution $`\mathrm{\Delta }P_{\mathrm{TT}}^{\mathrm{MH}}`$ remains finite.
## Chapter 5 Summary and Outlook
In this thesis, we computed the phase and the modulus of a composite, adjoint Higgs field $`\varphi `$ relevant for the thermodynamical description of an SU(2) Yang-Mills theory. The phase was defined by giving a spatial and moduli space average over an adjointly transforming two-point function and demanding BPS-saturation. The modulus could be inferred after assuming the existence of an externally given scale. It was seen that the field $`\varphi `$ exploits the instantaneous long range correlations in the classical caloron solution. The nontrivial temporal winding was obtained only after averaging over the entire admissible part of the moduli space. From the BPS equation, $`\varphi `$’s potential was uniquely deduced.
The two-loop contribution to the pressure of SU(2) Yang-Mills theory was computed. There are four contributing Feynman diagrams, among them one nonlocal. This one is seen to be dominant. The two-loop corrections are smaller than $`210^3`$ as compared to the one-loop result. Therefore the underlying picture of only very weakly interacting thermal quasiparticles is confirmed. The importance of the compositeness constraints in this process was pointed out. A microscopic interpretation in terms of TLM modes scattering off magnetic monopoles was given.
We saw that the compositeness constraints rule out (at least in the large coupling regime) a number of diagrams which one would naively expect to contribute. It would be interesting to consider this for higher loop-diagrams. Possibly, from some loop order upward only a few classes of diagrams will survive. The effects of the two-loop corrections for the pressure on the evolution of the gauge coupling are currently being worked on . A review on the discussed approach to SU(2) and SU(3) Yang-Mills thermodynamics (containing also the material presented here) and especially indicating implications on particle physics and cosmology is available in .
## Appendix A Appendix to Chapter 3
### A.1 Notation and conventions
In Chapter 3, the following conventions are used: SU(2) gauge field and field strength are written in matrix notation as
$$A_\mu =A_\mu ^a\frac{\lambda ^a}{2}\text{and}F_{\mu \nu }=F_{\mu \nu }^a\frac{\lambda ^a}{2},$$
(A.1)
where $`\lambda ^a`$ are the Pauli matrices. The covariant derivative is defined as
$$D_\mu =_\mu iA_\mu ,$$
(A.2)
where the coupling constant is absorbed in the gauge field. The field strength is in matrix notation
$$F_{\mu \nu }=_\mu A_\nu _\nu A_\mu i[A_\mu ,A_\nu ]$$
(A.3)
and in components
$$F_{\mu \nu }^a=_\mu A_\nu ^a_\nu A_\mu ^a+\epsilon ^{abc}A_\mu ^bA_\nu ^c.$$
(A.4)
We work in a Euclidean spacetime with $`x=(\tau ,𝒙)`$. Latin indices run from 1 to 3, Greek indices from 1 to 4. Upper and lower indices are not distinguished. The summation convention is implied if not specified otherwise. The totally antisymmetric symbol $`\epsilon _{\mu \nu \rho \sigma }`$ has $`\epsilon _{1234}=\epsilon ^{1234}=+1`$. The four-gradient is
$$_\mu =\frac{}{x_\mu },$$
(A.5)
derivatives with respect to time and space coordinates are written as
$$_4=\frac{}{\tau }\text{and}_i=\frac{}{x_i}$$
(A.6)
respectively. The derivative with respect to the radial coordinate $`r=\left|𝒙\right|`$ is written as
$$_r=\frac{}{r}.$$
(A.7)
The Pauli matrices are denoted by $`\lambda ^a`$, they have the explicit form
$$\lambda ^1=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right),\lambda ^2=\left(\begin{array}{cc}0& i\\ i& 0\end{array}\right),\lambda ^3=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right).$$
(A.8)
The Pauli matrices fulfill
$$\begin{array}{cc}\hfill [\lambda ^a,\lambda ^b]& =2i\epsilon _{abc}\lambda ^c\hfill \\ \hfill \lambda ^a\lambda ^b& =\delta _{ab}+i\epsilon _{abc}\lambda ^c\hfill \\ \hfill \lambda ^a\lambda ^b\lambda ^c& =\delta _{ab}\lambda ^c\delta _{ac}\lambda ^b+\delta _{bc}\lambda ^a+i\epsilon _{abc}.\hfill \end{array}$$
(A.9)
The traces of products of Pauli-matrices are
$$\begin{array}{cc}\hfill \mathrm{tr}\lambda ^a& =0\hfill \\ \hfill \mathrm{tr}\lambda ^a\lambda ^b& =2\delta _{ab}\hfill \\ \hfill \mathrm{tr}\lambda ^a\lambda ^b\lambda ^c& =2i\epsilon _{abc}\hfill \\ \hfill \mathrm{tr}\lambda ^a\lambda ^b\lambda ^c\lambda ^d& =2\left(\delta _{ab}\delta _{cd}+\delta _{ad}\delta _{bc}\delta _{ac}\delta _{bd}\right)\hfill \\ \hfill \mathrm{tr}\lambda ^a\lambda ^b\lambda ^c\lambda ^d\lambda ^e& =2i\left(\delta _{ab}\epsilon _{cde}+\delta _{cd}\epsilon _{abe}\delta _{cd}\epsilon _{abd}+\delta _{de}\epsilon _{abc}\right).\hfill \end{array}$$
(A.10)
The ’t Hooft symbols $`\eta `$ and $`\overline{\eta }`$ are defined as
$$\begin{array}{cc}\hfill \eta _{\mu \nu }^a& =\epsilon _{a\mu \nu }+\delta _{a\mu }\delta _{\nu 4}\delta _{a\nu }\delta _{\mu 4}\hfill \\ \hfill \overline{\eta }_{\mu \nu }^a& =\epsilon _{a\mu \nu }\delta _{a\mu }\delta _{\nu 4}+\delta _{a\nu }\delta _{\mu 4}.\hfill \end{array}$$
(A.11)
The symbols $`\eta `$ ($`\overline{\eta }`$) are (anti) self-dual and antisymmetric in the vector indices,
$$\begin{array}{cc}\hfill \eta _{\mu \nu }^a& =\frac{1}{2}\epsilon _{\mu \nu \alpha \beta }\eta _{\alpha \beta }^a\hfill \\ \hfill \overline{\eta }_{\mu \nu }^a& =\frac{1}{2}\epsilon _{\mu \nu \alpha \beta }\overline{\eta }_{\alpha \beta }^a\hfill \\ \hfill \eta _{\mu \nu }^a& =\eta _{\nu \mu }^a\hfill \\ \hfill \overline{\eta }_{\mu \nu }^a& =\overline{\eta }_{\nu \mu }^a.\hfill \end{array}$$
(A.12)
They fulfill a number of useful relations:
$$\begin{array}{cc}\hfill \eta _{\mu \nu }^a\eta _{\mu \nu }^b& =4\delta _{ab}\hfill \\ \hfill \eta _{\lambda \mu }^a\eta _{\lambda \nu }^a& =3\delta _{\mu \nu }\hfill \\ \hfill \eta _{\mu \nu }^a\eta _{\mu \nu }^a& =12\hfill \\ \hfill \eta _{\lambda \mu }^a\eta _{\lambda \nu }^b& =\delta _{ab}\delta _{\mu \nu }+\epsilon _{abc}\eta _{\mu \nu }^c\hfill \\ \hfill \epsilon _{abc}\eta _{\mu \nu }^b\eta _{\kappa \lambda }^c& =\delta _{\mu \kappa }\eta _{\nu \lambda }^a+\delta _{\nu \lambda }\eta _{\mu \kappa }^a\delta _{\mu \lambda }\eta _{\nu \kappa }^a\delta _{\nu \kappa }\eta _{\mu \lambda }^a\hfill \\ \hfill \epsilon _{abc}\eta _{\mu \nu }^b\eta _{\mu \lambda }^c& =2\eta _{\nu \lambda }^a\hfill \\ \hfill \eta _{\mu \nu }^a\eta _{\nu \lambda }^b\eta _{\lambda \mu }^c& =4\epsilon _{abc}.\hfill \end{array}$$
(A.13)
The relations (A.13) hold for $`\overline{\eta }`$ as well. Besides that,
$$\begin{array}{cc}\hfill \eta _{\mu \nu }^a\eta _{\kappa \lambda }^a& =\delta _{\mu \kappa }\delta _{\nu \lambda }\delta _{\mu \lambda }\delta _{\nu \kappa }+\epsilon _{\mu \nu \kappa \lambda }\hfill \\ \hfill \overline{\eta }_{\mu \nu }^a\overline{\eta }_{\kappa \lambda }^a& =\delta _{\mu \kappa }\delta _{\nu \lambda }\delta _{\mu \lambda }\delta _{\nu \kappa }\epsilon _{\mu \nu \kappa \lambda }\hfill \\ \hfill \epsilon _{\lambda \mu \nu \sigma }\eta _{\rho \sigma }^a& =\delta _{\rho \lambda }\eta _{\mu \nu }^a+\delta _{\rho \nu }\eta _{\lambda \mu }^a+\delta _{\rho \mu }\eta _{\nu \kappa }^a\hfill \\ \hfill \epsilon _{\lambda \mu \nu \sigma }\overline{\eta }_{\rho \sigma }^a& =\delta _{\rho \lambda }\overline{\eta }_{\mu \nu }^a+\delta _{\rho \nu }\overline{\eta }_{\lambda \mu }^a+\delta _{\rho \mu }\overline{\eta }_{\nu \kappa }^a\hfill \\ \hfill \eta _{\mu \nu }^a\overline{\eta }_{\mu \nu }^b& =0.\hfill \end{array}$$
(A.14)
### A.2 Caloron field strength
The SU(2) field strength in the ’nonperturbative’ convention is
$$F_{\mu \nu }^a=_\mu A_\nu ^a_\nu A_\mu ^a+\epsilon ^{abc}A_\mu ^bA_\nu ^c.$$
(A.15)
Inserting the singular gauge instanton field
$$A_\mu ^a=\overline{\eta }_{\mu \kappa }^a_\kappa \mathrm{ln}\mathrm{\Pi }=\overline{\eta }_{\mu \kappa }^a\frac{_\kappa \mathrm{\Pi }}{\mathrm{\Pi }},$$
(A.16)
we get
$`F_{\mu \nu }^a`$ $`=\overline{\eta }_{\nu \kappa }^a_\mu {\displaystyle \frac{_\kappa \mathrm{\Pi }}{\mathrm{\Pi }}}+\overline{\eta }_{\mu \kappa }^a_\nu {\displaystyle \frac{_\kappa \mathrm{\Pi }}{\mathrm{\Pi }}}+\epsilon ^{abc}\overline{\eta }_{\mu \kappa }^b\overline{\eta }_{\nu \lambda }^c{\displaystyle \frac{\left(_\kappa \mathrm{\Pi }\right)\left(_\lambda \mathrm{\Pi }\right)}{\mathrm{\Pi }^2}}`$
$`=\overline{\eta }_{\nu \kappa }^a_\mu {\displaystyle \frac{_\kappa \mathrm{\Pi }}{\mathrm{\Pi }}}+\overline{\eta }_{\mu \kappa }^a_\nu {\displaystyle \frac{_\kappa \mathrm{\Pi }}{\mathrm{\Pi }}}+\left(\overline{\eta }_{\mu \nu }^a\delta _{\kappa \lambda }+\overline{\eta }_{\kappa \lambda }^a\delta _{\mu \nu }\overline{\eta }_{\mu \lambda }^a\delta _{\kappa \nu }\overline{\eta }_{\kappa \nu }^a\delta _{\mu \lambda }\right){\displaystyle \frac{\left(_\kappa \mathrm{\Pi }\right)\left(_\lambda \mathrm{\Pi }\right)}{\mathrm{\Pi }^2}}`$
$`=\overline{\eta }_{\nu \kappa }^a{\displaystyle \frac{\mathrm{\Pi }\left(_\mu _\kappa \mathrm{\Pi }\right)\left(_\kappa \mathrm{\Pi }\right)\left(_\mu \mathrm{\Pi }\right)}{\mathrm{\Pi }^2}}+\overline{\eta }_{\mu \kappa }^a{\displaystyle \frac{\mathrm{\Pi }\left(_\nu _\kappa \mathrm{\Pi }\right)\left(_\kappa \mathrm{\Pi }\right)\left(_\nu \mathrm{\Pi }\right)}{\mathrm{\Pi }^2}}`$ (A.17)
$`+\left(\overline{\eta }_{\mu \nu }^a{\displaystyle \frac{\left(_\kappa \mathrm{\Pi }\right)\left(_\kappa \mathrm{\Pi }\right)}{\mathrm{\Pi }^2}}\overline{\eta }_{\mu \lambda }^a{\displaystyle \frac{\left(_\nu \mathrm{\Pi }\right)\left(_\lambda \mathrm{\Pi }\right)}{\mathrm{\Pi }^2}}\overline{\eta }_{\kappa \nu }^a{\displaystyle \frac{\left(_\kappa \mathrm{\Pi }\right)\left(_\mu \mathrm{\Pi }\right)}{\mathrm{\Pi }^2}}\right)`$
$`=\overline{\eta }_{\mu \nu }^a{\displaystyle \frac{\left(_\kappa \mathrm{\Pi }\right)\left(_\kappa \mathrm{\Pi }\right)}{\mathrm{\Pi }^2}}\overline{\eta }_{\nu \kappa }^a{\displaystyle \frac{\mathrm{\Pi }\left(_\mu _\kappa \mathrm{\Pi }\right)2\left(_\kappa \mathrm{\Pi }\right)\left(_\mu \mathrm{\Pi }\right)}{\mathrm{\Pi }^2}}+\overline{\eta }_{\mu \kappa }^a{\displaystyle \frac{\mathrm{\Pi }\left(_\nu _\kappa \mathrm{\Pi }\right)2\left(_\kappa \mathrm{\Pi }\right)\left(_\nu \mathrm{\Pi }\right)}{\mathrm{\Pi }^2}}`$
For convenience, we define the quantities
$$P=\frac{\left(_\kappa \mathrm{\Pi }\right)\left(_\kappa \mathrm{\Pi }\right)}{\mathrm{\Pi }^2}\text{and}P_{\mu \nu }=\frac{\mathrm{\Pi }\left(_\mu _\nu \mathrm{\Pi }\right)2\left(_\mu \mathrm{\Pi }\right)\left(_\nu \mathrm{\Pi }\right)}{\mathrm{\Pi }^2},$$
(A.18)
such that the instanton field strength reads
$$F_{\mu \nu }^a=\overline{\eta }_{\mu \nu }^aP\overline{\eta }_{\nu \kappa }^aP_{\mu \kappa }+\overline{\eta }_{\mu \kappa }^aP_{\nu \kappa }.$$
(A.19)
The symbols (A.18) have the properties
$$P_{\mu \nu }=P_{\nu \mu }\text{and}P_{\mu \mu }=2P.$$
(A.20)
The pre-potential $`\mathrm{\Pi }`$ for a single caloron is given in Eq. (2.104). Using the fact that it depends on the radial coordinate $`r=\left|𝒙\right|`$ only (and on Euclidean time), $`\mathrm{\Pi }(\tau ,𝒙)=\mathrm{\Pi }(\tau ,r)`$, $`P`$ and $`P_{\mu \nu }`$ can be expressed in terms of the derivatives of $`\mathrm{\Pi }`$ with respect to $`\tau `$ and $`r`$. In particular, the Laplace equation
$$\mathrm{\Pi }^1_\mu _\mu \mathrm{\Pi }=0,$$
(A.21)
which results from the demand for self-duality, translates into
$$\frac{_4^2\mathrm{\Pi }}{\mathrm{\Pi }}+\frac{_r^2\mathrm{\Pi }}{\mathrm{\Pi }}+2\frac{_r\mathrm{\Pi }}{r\mathrm{\Pi }}=0.$$
(A.22)
At an arbitrary point $`x=(\tau ,𝒙)`$, the components of $`P_{\mu \nu }`$ are explicitly given as
$$\begin{array}{cc}\hfill P_{44}\left(x\right)& =\frac{_4^2\mathrm{\Pi }\left(x\right)}{\mathrm{\Pi }\left(x\right)}2\frac{\left[_4\mathrm{\Pi }\left(x\right)\right]^2}{\mathrm{\Pi }^2\left(x\right)}\hfill \\ \hfill P_{4i}\left(x\right)& =\frac{x_i}{r}\left(\frac{_4_r\mathrm{\Pi }\left(x\right)}{\mathrm{\Pi }\left(x\right)}2\frac{\left[_4\mathrm{\Pi }\left(x\right)\right]\left[_r\mathrm{\Pi }\left(x\right)\right]}{\mathrm{\Pi }^2\left(x\right)}\right)\hfill \\ \hfill P_{ij}\left(x\right)& =\delta _{ij}\frac{_r\mathrm{\Pi }\left(x\right)}{r\mathrm{\Pi }\left(x\right)}+\frac{x_ix_j}{r^2}\left(\frac{_r^2\mathrm{\Pi }\left(x\right)}{\mathrm{\Pi }\left(x\right)}\frac{_r\mathrm{\Pi }\left(x\right)}{r\mathrm{\Pi }\left(x\right)}2\frac{\left[_r\mathrm{\Pi }\left(x\right)\right]^2}{\mathrm{\Pi }^2\left(x\right)}\right),\hfill \end{array}$$
(A.23)
and
$$P\left(x\right)=\frac{\left[_4\mathrm{\Pi }\left(x\right)\right]^2}{\mathrm{\Pi }^2\left(x\right)}+\frac{\left[_r\mathrm{\Pi }\left(x\right)\right]^2}{\mathrm{\Pi }^2\left(x\right)}.$$
(A.24)
At $`x=(\tau ,0)`$, the tensor $`P_{\mu \nu }(\tau ,0)`$ is diagonal and its components are
$$\begin{array}{cc}\hfill P_{4i}(\tau ,0)& =P_{i4}(\tau ,0)=0\hfill \\ \hfill P_{44}(\tau ,0)& =\frac{_4^2\mathrm{\Pi }(\tau ,0)}{\mathrm{\Pi }(\tau ,0)}2\frac{\left[_4\mathrm{\Pi }(\tau ,0)\right]^2}{\mathrm{\Pi }^2(\tau ,0)}\hfill \\ \hfill P_{ij}(\tau ,0)& =\frac{1}{3}\delta _{ij}\frac{_4^2\mathrm{\Pi }(\tau ,0)}{\mathrm{\Pi }(\tau ,0)},\hfill \end{array}$$
(A.25)
and
$$P(\tau ,0)=\frac{\left[_4\mathrm{\Pi }(\tau ,0)\right]^2}{\mathrm{\Pi }^2(\tau ,0)}.$$
(A.26)
### A.3 Details to section 3.2.3
All spacetime dependent objects are to be evaluated at the same time $`\tau `$, so that we may, for simplicity of notation, suppress the time dependence and write $`F_{\mu \nu }\left(𝒙\right)`$ instead of $`F_{\mu \nu }(\tau ,𝒙)`$ and $`\{0,𝒙\}`$ instead of $`\{(\tau ,0),(\tau ,𝒙)\}`$ etc. Besides that, we abbreviate the time derivative with a dot and the radial derivative with a prime.
In the following some expressions containing $`P_{\mu \nu }`$ and $`P`$, which we will need later, are calculated:
$$\begin{array}{cc}\hfill \overline{\eta }_{\mu \kappa }^b\overline{\eta }_{\nu \rho }^bP_{\nu \kappa }\left(0\right)P_{\mu \rho }\left(𝒙\right)& =\left(\delta _{\mu \nu }\delta _{\kappa \rho }\delta _{\mu \rho }\delta _{\nu \kappa }\right)P_{\nu \kappa }\left(0\right)P_{\mu \rho }\left(𝒙\right)\hfill \\ & =P_{\mu \nu }\left(0\right)P_{\mu \nu }\left(𝒙\right)P_{\nu \nu }\left(0\right)P_{\mu \mu }\left(𝒙\right)\hfill \\ & =P_{\mu \nu }\left(0\right)P_{\mu \nu }\left(𝒙\right)4P\left(0\right)P\left(𝒙\right)\hfill \end{array}$$
(A.27)
$$\begin{array}{cc}\hfill \overline{\eta }_{\mu \kappa }^a\overline{\eta }_{\nu \rho }^cP_{\nu \kappa }\left(0\right)P_{\mu \rho }\left(𝒙\right)& =\overline{\eta }_{\mu 4}^a\overline{\eta }_{4\rho }^cP_{44}\left(0\right)P_{\mu \rho }\left(𝒙\right)+\overline{\eta }_{\mu k}^a\overline{\eta }_{n\rho }^cP_{nk}\left(0\right)P_{\mu \rho }\left(𝒙\right)\hfill \\ & =P_{44}\left(0\right)P_{ac}\left(𝒙\right)\frac{1}{3}\overline{\eta }_{\mu n}^a\overline{\eta }_{n\rho }^c\frac{\ddot{\mathrm{\Pi }}\left(0\right)}{\mathrm{\Pi }\left(0\right)}P_{\mu \rho }\left(𝒙\right)\hfill \\ & =P_{44}\left(0\right)P_{ac}\left(𝒙\right)\frac{1}{3}\left(\epsilon _{a\mu n}+\delta _{\mu 4}\delta _{an}\right)\left(\epsilon _{cn\rho }\delta _{cn}\delta _{\rho 4}\right)\frac{\ddot{\mathrm{\Pi }}\left(0\right)}{\mathrm{\Pi }\left(0\right)}P_{\mu \rho }\left(𝒙\right)\hfill \\ & =P_{44}\left(0\right)P_{ac}\left(𝒙\right)\frac{1}{3}\frac{\ddot{\mathrm{\Pi }}\left(0\right)}{\mathrm{\Pi }\left(0\right)}\left(P_{ac}\left(𝒙\right)\delta _{ac}P_{ii}\left(𝒙\right)\delta _{ac}P_{44}\left(𝒙\right)\right)\hfill \\ & =P_{44}\left(0\right)P_{ac}\left(𝒙\right)\frac{1}{3}\frac{\ddot{\mathrm{\Pi }}\left(0\right)}{\mathrm{\Pi }\left(0\right)}\left(P_{ac}\left(𝒙\right)\delta _{ac}P_{\mu \mu }\left(𝒙\right)\right)\hfill \\ & =P_{44}\left(0\right)P_{ac}\left(𝒙\right)\frac{1}{3}\left(P_{44}\left(0\right)+2P\left(0\right)\right)\left(P_{ac}\left(𝒙\right)+2\delta _{ac}P\left(𝒙\right)\right)\hfill \\ & =\frac{2}{3}\left(2P_{44}\left(0\right)+P\left(0\right)\right)P_{ac}\left(𝒙\right)\frac{2}{3}\left(P_{44}\left(0\right)+2P\left(0\right)\right)\delta _{ac}P\left(𝒙\right)\hfill \end{array}$$
(A.28)
$$P_{ii}\left(𝒙\right)=\frac{\ddot{\mathrm{\Pi }}\left(r\right)}{\mathrm{\Pi }\left(r\right)}2\frac{\mathrm{\Pi }^2\left(r\right)}{\mathrm{\Pi }^2\left(r\right)}$$
(A.29)
$$\begin{array}{cc}\hfill P_{\mu \nu }\left(0\right)P_{\mu \nu }\left(𝒙\right)& =P_{44}\left(0\right)P_{44}\left(𝒙\right)+P_{mn}\left(0\right)P_{mn}\left(𝒙\right)\hfill \\ & =P_{44}\left(0\right)P_{44}\left(𝒙\right)\frac{1}{3}\frac{\ddot{\mathrm{\Pi }}\left(0\right)}{\mathrm{\Pi }\left(0\right)}P_{ii}\left(𝒙\right)\hfill \\ & =\left(\frac{\ddot{\mathrm{\Pi }}\left(0\right)}{\mathrm{\Pi }\left(0\right)}2\frac{\dot{\mathrm{\Pi }}^2\left(0\right)}{\mathrm{\Pi }^2\left(0\right)}\right)\left(\frac{\ddot{\mathrm{\Pi }}\left(r\right)}{\mathrm{\Pi }\left(r\right)}2\frac{\dot{\mathrm{\Pi }}^2\left(r\right)}{\mathrm{\Pi }^2\left(r\right)}\right)+\frac{1}{3}\frac{\ddot{\mathrm{\Pi }}\left(0\right)}{\mathrm{\Pi }\left(0\right)}\left(\frac{\ddot{\mathrm{\Pi }}\left(r\right)}{\mathrm{\Pi }\left(r\right)}+2\frac{\mathrm{\Pi }^2\left(r\right)}{\mathrm{\Pi }^2\left(r\right)}\right)\hfill \end{array}$$
(A.30)
$$P_{\mu \nu }\left(0\right)P_{\mu \nu }\left(𝒙\right)P\left(0\right)P\left(𝒙\right)=\left(\frac{\dot{\mathrm{\Pi }}^2\left(0\right)}{\mathrm{\Pi }^2\left(0\right)}\frac{2}{3}\frac{\ddot{\mathrm{\Pi }}\left(0\right)}{\mathrm{\Pi }\left(0\right)}\right)\left(3\frac{\dot{\mathrm{\Pi }}^2\left(r\right)}{\mathrm{\Pi }^2\left(r\right)}2\frac{\ddot{\mathrm{\Pi }}\left(r\right)}{\mathrm{\Pi }\left(r\right)}\frac{\mathrm{\Pi }^2\left(r\right)}{\mathrm{\Pi }^2\left(r\right)}\right)$$
(A.31)
$$\begin{array}{cc}\hfill \overline{\eta }_{\kappa \rho }^aP_{\mu \kappa }\left(0\right)P_{\mu \rho }\left(𝒙\right)& =\overline{\eta }_{4\rho }^aP_{44}\left(0\right)P_{4\rho }\left(𝒙\right)+\overline{\eta }_{k\rho }^aP_{mk}\left(0\right)P_{m\rho }\left(𝒙\right)\hfill \\ & =P_{44}\left(0\right)P_{4a}\left(𝒙\right)\frac{1}{3}\frac{\ddot{\mathrm{\Pi }}\left(0\right)}{\mathrm{\Pi }\left(0\right)}\overline{\eta }_{m\rho }^aP_{m\rho }\left(𝒙\right)\hfill \\ & =P_{44}\left(0\right)P_{4a}\left(𝒙\right)+\frac{1}{3}\frac{\ddot{\mathrm{\Pi }}\left(0\right)}{\mathrm{\Pi }\left(0\right)}P_{4a}\left(𝒙\right)\hfill \\ & =\left(P_{44}\left(0\right)+\frac{1}{3}\frac{\ddot{\mathrm{\Pi }}\left(0\right)}{\mathrm{\Pi }\left(0\right)}\right)P_{4a}\left(𝒙\right)\hfill \\ & =2\frac{x^a}{r}\left(\frac{\dot{\mathrm{\Pi }}^{}\left(r\right)}{\mathrm{\Pi }\left(r\right)}2\frac{\dot{\mathrm{\Pi }}\left(r\right)\mathrm{\Pi }^{}\left(r\right)}{\mathrm{\Pi }^2\left(r\right)}\right)\left(\frac{2}{3}\frac{\ddot{\mathrm{\Pi }}\left(0\right)}{\mathrm{\Pi }\left(0\right)}2\frac{\dot{\mathrm{\Pi }}^2\left(0\right)}{\mathrm{\Pi }^2\left(0\right)}\right)\hfill \end{array}$$
(A.32)
$$\begin{array}{cc}\hfill \frac{x^ax^f}{r^2}\overline{\eta }_{\kappa \rho }^fP_{\mu \kappa }\left(0\right)P_{\mu \rho }\left(𝒙\right)& =2\frac{x^a}{r}\left(\frac{\dot{\mathrm{\Pi }}^{}\left(r\right)}{\mathrm{\Pi }\left(r\right)}2\frac{\dot{\mathrm{\Pi }}\left(r\right)\mathrm{\Pi }^{}\left(r\right)}{\mathrm{\Pi }^2\left(r\right)}\right)\left(\frac{2}{3}\frac{\ddot{\mathrm{\Pi }}\left(0\right)}{\mathrm{\Pi }\left(0\right)}2\frac{\dot{\mathrm{\Pi }}^2\left(0\right)}{\mathrm{\Pi }^2\left(0\right)}\right)\hfill \end{array}$$
(A.33)
Integrand of (3.31). Inserting the result for the Wilson lines from Eq. (3.37) and writing the field strength in components, we have Eq. (3.42), which in simplified notation reads
$`\mathrm{tr}\lambda ^aF_{\mu \nu }\left(0\right)\{0,𝒙\}F_{\mu \nu }\left(𝒙\right)\{𝒙,0\}`$
$`={\displaystyle \frac{1}{2}}\mathrm{tr}\left[\lambda ^a\lambda ^b\left(\mathrm{cos}g\left(r\right)+i\lambda ^c{\displaystyle \frac{x^c}{r}}\mathrm{sin}g\left(r\right)\right)\lambda ^d\left(\mathrm{cos}g\left(r\right)i\lambda ^e{\displaystyle \frac{x^e}{r}}\mathrm{sin}g\left(r\right)\right)\right]F_{\mu \nu }^b\left(0\right)F_{\mu \nu }^d\left(𝒙\right)`$
$`={\displaystyle \frac{1}{2}}\mathrm{cos}^2g\left(r\right)\mathrm{tr}\left[\lambda ^a\lambda ^b\lambda ^d\right]F_{\mu \nu }^b\left(0\right)F_{\mu \nu }^d\left(𝒙\right)`$
$`{\displaystyle \frac{i}{2}}\mathrm{sin}g\left(r\right)\mathrm{cos}g\left(r\right)\mathrm{tr}\left[\lambda ^a\lambda ^b\lambda ^d\lambda ^e\right]{\displaystyle \frac{x^e}{r}}F_{\mu \nu }^b\left(0\right)F_{\mu \nu }^d\left(𝒙\right)`$ (A.34)
$`+{\displaystyle \frac{i}{2}}\mathrm{sin}g\left(r\right)\mathrm{cos}g\left(r\right)\mathrm{tr}\left[\lambda ^a\lambda ^b\lambda ^c\lambda ^d\right]{\displaystyle \frac{x^c}{r}}F_{\mu \nu }^b\left(0\right)F_{\mu \nu }^d\left(𝒙\right)`$
$`+{\displaystyle \frac{1}{2}}\mathrm{sin}^2g\left(r\right)\mathrm{tr}\left[\lambda ^a\lambda ^b\lambda ^c\lambda ^d\lambda ^e\right]{\displaystyle \frac{x^cx^e}{r^2}}F_{\mu \nu }^b\left(0\right)F_{\mu \nu }^d\left(𝒙\right).`$
We will calculate each of the four summands separately. Therefore we write the Lorentz-trace of the field components $`F_{\mu \nu }^b\left(0\right)F_{\mu \nu }^d\left(𝒙\right)`$ in terms of the symbols $`P`$ and $`P_{\mu \nu }`$ introduced in Eqs. (A.18),
$`F_{\mu \nu }^b\left(0\right)F_{\mu \nu }^d\left(𝒙\right)`$
$`=\left[\overline{\eta }_{\mu \nu }^bP\left(0\right)+\overline{\eta }_{\mu \kappa }^bP_{\nu \kappa }\left(0\right)\overline{\eta }_{\nu \kappa }^bP_{\mu \kappa }\left(0\right)\right]\left[\overline{\eta }_{\mu \nu }^dP\left(𝒙\right)+\overline{\eta }_{\mu \rho }^dP_{\nu \rho }\left(𝒙\right)\overline{\eta }_{\nu \rho }^dP_{\mu \rho }\left(𝒙\right)\right]`$ (A.35)
$`=2\left\{\delta _{bd}\left[P_{\mu \nu }\left(0\right)P_{\mu \nu }\left(𝒙\right)2P\left(0\right)P\left(𝒙\right)\right]+\epsilon _{bdf}\overline{\eta }_{\kappa \rho }^fP_{\mu \kappa }\left(0\right)P_{\mu \rho }\left(𝒙\right)\overline{\eta }_{\mu \kappa }^b\overline{\eta }_{\nu \rho }^dP_{\nu \kappa }\left(0\right)P_{\mu \rho }\left(𝒙\right)\right\}.`$
Furthermore, we need the traces of products of Pauli matrices Eqs. (A.10), some of the relations for ’t Hooft symbols from Sec. A.1, and the expressions prepared in Eqs. (A.27) through (A.33). The first of the four terms in Eq. (A.3) is
$`{\displaystyle \frac{1}{2}}\mathrm{cos}^2g\left(r\right)\mathrm{tr}\left[\lambda ^a\lambda ^b\lambda ^d\right]F_{\mu \nu }^b\left(0\right)F_{\mu \nu }^d\left(𝒙\right)`$
$`=i\mathrm{cos}^2g\left(r\right)\epsilon ^{abd}F_{\mu \nu }^b\left(0\right)F_{\mu \nu }^d\left(𝒙\right)`$
$`=2i\mathrm{cos}^2g\left(r\right)\overline{\eta }_{\kappa \rho }^aP_{\mu \kappa }\left(0\right)P_{\mu \rho }\left(𝒙\right)`$
$`=4i\mathrm{cos}^2g\left(r\right){\displaystyle \frac{x^a}{r}}\left({\displaystyle \frac{\dot{\mathrm{\Pi }}^2\left(0\right)}{\mathrm{\Pi }^2\left(0\right)}}{\displaystyle \frac{2}{3}}{\displaystyle \frac{\ddot{\mathrm{\Pi }}\left(0\right)}{\mathrm{\Pi }\left(0\right)}}\right)\left({\displaystyle \frac{\dot{\mathrm{\Pi }}^{}\left(r\right)}{\mathrm{\Pi }\left(r\right)}}2{\displaystyle \frac{\dot{\mathrm{\Pi }}\left(r\right)\mathrm{\Pi }^{}\left(r\right)}{\mathrm{\Pi }^2\left(r\right)}}\right).`$ (A.36)
In the third line, we used that
$$\epsilon ^{abd}\overline{\eta }_{\mu \kappa }^b\overline{\eta }_{\nu \rho }^dP_{\nu \kappa }\left(0\right)P_{\mu \rho }\left(𝒙\right)=\epsilon ^{adb}\overline{\eta }_{\rho \nu }^d\overline{\eta }_{\kappa \mu }^bP_{\kappa \nu }\left(0\right)P_{\rho \mu }\left(𝒙\right)=\epsilon ^{abd}\overline{\eta }_{\mu \kappa }^b\overline{\eta }_{\nu \rho }^dP_{\nu \kappa }\left(0\right)P_{\mu \rho }\left(𝒙\right)$$
vanishes. The second summand is
$`{\displaystyle \frac{i}{2}}\mathrm{cos}g\left(r\right)\mathrm{sin}g\left(r\right)\mathrm{tr}\left[\lambda ^a\lambda ^b\lambda ^d\lambda ^e\right]{\displaystyle \frac{x^e}{r}}F_{\mu \nu }^b\left(0\right)F_{\mu \nu }^d\left(𝒙\right)`$
$`=i\mathrm{cos}g\left(r\right)\mathrm{sin}g\left(r\right)\left(\delta ^{ab}\delta ^{de}\delta ^{ac}\delta ^{be}+\delta ^{bd}\delta ^{ae}\right){\displaystyle \frac{x^e}{r}}F_{\mu \nu }^b\left(0\right)F_{\mu \nu }^d\left(𝒙\right)`$
$`=i\mathrm{cos}g\left(r\right)\mathrm{sin}g\left(r\right){\displaystyle \frac{x^e}{r}}\{3\delta _{ae}[P_{\mu \nu }\left(0\right)P_{\mu \nu }\left(𝒙\right)2P\left(0\right)P\left(𝒙\right)]+2\epsilon _{aef}\overline{\eta }_{\kappa \rho }^fP_{\mu \kappa }\left(0\right)P_{\mu \rho }\left(𝒙\right)`$
$`(\overline{\eta }_{\mu \kappa }^a\overline{\eta }_{\nu \rho }^e\overline{\eta }_{\mu \kappa }^e\overline{\eta }_{\nu \rho }^a)P_{\nu \kappa }\left(0\right)P_{\mu \rho }\left(𝒙\right)\delta _{ae}\overline{\eta }_{\mu \kappa }^b\overline{\eta }_{\nu \rho }^bP_{\nu \kappa }\left(0\right)P_{\mu \rho }\left(𝒙\right)\}`$
$`=i\mathrm{cos}g\left(r\right)\mathrm{sin}g\left(r\right){\displaystyle \frac{x^a}{r}}\left\{3\left[P_{\mu \nu }\left(0\right)P_{\mu \nu }\left(𝒙\right)2P\left(0\right)P\left(𝒙\right)\right]\left[P_{\mu \nu }\left(0\right)P_{\mu \nu }\left(𝒙\right)4P\left(0\right)P\left(𝒙\right)\right]\right\}`$
$`=2i\mathrm{cos}g\left(r\right)\mathrm{sin}g\left(r\right){\displaystyle \frac{x^a}{r}}\left\{P_{\mu \nu }\left(0\right)P_{\mu \nu }\left(𝒙\right)P\left(0\right)P\left(𝒙\right)\right\}`$
$`=2i\mathrm{cos}g\left(r\right)\mathrm{sin}g\left(r\right){\displaystyle \frac{x^a}{r}}\left({\displaystyle \frac{\dot{\mathrm{\Pi }}^2\left(0\right)}{\mathrm{\Pi }^2\left(0\right)}}{\displaystyle \frac{2}{3}}{\displaystyle \frac{\ddot{\mathrm{\Pi }}\left(0\right)}{\mathrm{\Pi }\left(0\right)}}\right)\left(3{\displaystyle \frac{\dot{\mathrm{\Pi }}^2\left(r\right)}{\mathrm{\Pi }^2\left(r\right)}}2{\displaystyle \frac{\ddot{\mathrm{\Pi }}\left(r\right)}{\mathrm{\Pi }\left(r\right)}}{\displaystyle \frac{\mathrm{\Pi }^2\left(r\right)}{\mathrm{\Pi }^2\left(r\right)}}\right).`$ (A.37)
The third summand is
$`{\displaystyle \frac{i}{2}}\mathrm{cos}g\left(r\right)\mathrm{sin}g\left(r\right)\mathrm{tr}\left[\lambda ^a\lambda ^b\lambda ^c\lambda ^d\right]{\displaystyle \frac{x^c}{r}}F_{\mu \nu }^b\left(0\right)F_{\mu \nu }^d\left(𝒙\right)`$
$`=i\mathrm{cos}g\left(r\right)\mathrm{sin}g\left(r\right)\left(\delta ^{ab}\delta ^{cd}\delta ^{ac}\delta ^{bd}+\delta ^{bc}\delta ^{ad}\right){\displaystyle \frac{x^c}{r}}F_{\mu \nu }^b\left(0\right)F_{\mu \nu }^d\left(𝒙\right)`$
$`=i\mathrm{cos}g\left(r\right)\mathrm{sin}g\left(r\right){\displaystyle \frac{x^c}{r}}\{\delta _{ac}[P_{\mu \nu }\left(0\right)P_{\mu \nu }\left(𝒙\right)2P\left(0\right)P\left(𝒙\right)]`$
$`(\overline{\eta }_{\mu \kappa }^a\overline{\eta }_{\nu \rho }^c+\overline{\eta }_{\mu \kappa }^c\overline{\eta }_{\nu \rho }^a)P_{\nu \kappa }\left(0\right)P_{\mu \rho }\left(𝒙\right)+\delta _{ac}\overline{\eta }_{\mu \kappa }^b\overline{\eta }_{\nu \rho }^bP_{\nu \kappa }\left(0\right)P_{\mu \rho }\left(𝒙\right)\}`$
$`=2i\mathrm{cos}g\left(r\right)\mathrm{sin}g\left(r\right)\left\{{\displaystyle \frac{x^a}{r}}P\left(0\right)P\left(𝒙\right)+{\displaystyle \frac{x^c}{r}}\overline{\eta }_{\mu \kappa }^a\overline{\eta }_{\nu \rho }^cP_{\nu \kappa }\left(0\right)P_{\mu \rho }\left(𝒙\right)\right\}`$
$`=2i\mathrm{cos}g\left(r\right)\mathrm{sin}g\left(r\right){\displaystyle \frac{x^a}{r}}\left({\displaystyle \frac{\dot{\mathrm{\Pi }}^2\left(0\right)}{\mathrm{\Pi }^2\left(0\right)}}{\displaystyle \frac{2}{3}}{\displaystyle \frac{\ddot{\mathrm{\Pi }}\left(0\right)}{\mathrm{\Pi }\left(0\right)}}\right)\left({\displaystyle \frac{\dot{\mathrm{\Pi }}^2\left(r\right)}{\mathrm{\Pi }^2\left(r\right)}}3{\displaystyle \frac{\mathrm{\Pi }^2\left(r\right)}{\mathrm{\Pi }^2\left(r\right)}}+2{\displaystyle \frac{\mathrm{\Pi }^{\prime \prime }\left(r\right)}{\mathrm{\Pi }\left(r\right)}}\right).`$ (A.38)
The fourth and last one is
$`{\displaystyle \frac{1}{2}}\mathrm{sin}^2g\left(r\right)\mathrm{tr}\left[\lambda ^a\lambda ^b\lambda ^c\lambda ^d\lambda ^e\right]{\displaystyle \frac{x^cx^e}{r^2}}F_{\mu \nu }^b\left(0\right)F_{\mu \nu }^d\left(𝒙\right)`$
$`=i\mathrm{sin}^2g\left(r\right)\left(\delta _{ab}\epsilon _{cde}+\delta _{cd}\epsilon _{abe}\delta _{ce}\epsilon _{abd}+\delta _{de}\epsilon _{abc}\right){\displaystyle \frac{x^cx^e}{r^2}}F_{\mu \nu }^b\left(0\right)F_{\mu \nu }^d\left(𝒙\right)`$
$`=i\mathrm{sin}^2g\left(r\right)\left(2\epsilon _{abc}{\displaystyle \frac{x^cx^d}{r^2}}\epsilon _{abd}\right)F_{\mu \nu }^b\left(0\right)F_{\mu \nu }^d\left(𝒙\right)`$
$`=i\mathrm{sin}^2g\left(r\right)\left\{\left(2\epsilon _{abc}\epsilon _{bdf}{\displaystyle \frac{x^cx^d}{r^2}}2\delta _{af}\right)\overline{\eta }_{\kappa \rho }^fP_{\mu \kappa }\left(0\right)P_{\mu \rho }\left(𝒙\right)2\epsilon _{abc}{\displaystyle \frac{x^cx^d}{r^2}}\overline{\eta }_{\mu \kappa }^b\overline{\eta }_{\nu \rho }^dP_{\nu \kappa }\left(0\right)P_{\mu \rho }\left(𝒙\right)\right\}`$
$`=2i\mathrm{sin}^2g\left(r\right)\overline{\eta }_{\kappa \rho }^aP_{\mu \kappa }\left(0\right)P_{\mu \rho }\left(𝒙\right)`$
$`=4i\mathrm{sin}^2g\left(r\right){\displaystyle \frac{x^a}{r}}\left({\displaystyle \frac{\dot{\mathrm{\Pi }}^2\left(0\right)}{\mathrm{\Pi }^2\left(0\right)}}{\displaystyle \frac{2}{3}}{\displaystyle \frac{\ddot{\mathrm{\Pi }}\left(0\right)}{\mathrm{\Pi }\left(0\right)}}\right)\left({\displaystyle \frac{\dot{\mathrm{\Pi }}^{}\left(r\right)}{\mathrm{\Pi }\left(r\right)}}2{\displaystyle \frac{\dot{\mathrm{\Pi }}\left(r\right)\mathrm{\Pi }^{}\left(r\right)}{\mathrm{\Pi }^2\left(r\right)}}\right).`$ (A.39)
Summing Eqs. (A.3) through (A.3) and using
$$\mathrm{cos}^2\alpha \mathrm{sin}^2\alpha =\mathrm{cos}2\alpha \text{and}\mathrm{cos}\alpha \mathrm{sin}\alpha =\frac{1}{2}\mathrm{sin}2\alpha $$
(A.40)
yields (with the $`\tau `$’s reinserted)
$`\mathrm{tr}\lambda ^aF_{\mu \nu }(\tau ,0)\{(\tau ,0),(\tau ,𝒙)\}F_{\mu \nu }(\tau ,𝒙)\{(\tau ,𝒙),(\tau ,0)\}`$
$`=2i{\displaystyle \frac{x^a}{r}}({\displaystyle \frac{\dot{\mathrm{\Pi }}^2(\tau ,0)}{\mathrm{\Pi }^2(\tau ,0)}}{\displaystyle \frac{2}{3}}{\displaystyle \frac{\ddot{\mathrm{\Pi }}(\tau ,0)}{\mathrm{\Pi }(\tau ,0)}})\{2\mathrm{cos}\left(2g(\tau ,r)\right)(2{\displaystyle \frac{\dot{\mathrm{\Pi }}(\tau ,r)\mathrm{\Pi }^{}(\tau ,r)}{\mathrm{\Pi }^2(\tau ,r)}}{\displaystyle \frac{\dot{\mathrm{\Pi }}^{}(\tau ,r)}{\mathrm{\Pi }(\tau ,r)}})`$
$`+\mathrm{sin}\left(2g(\tau ,r)\right)(2{\displaystyle \frac{\mathrm{\Pi }^2(\tau ,r)}{\mathrm{\Pi }^2(\tau ,r)}}2{\displaystyle \frac{\dot{\mathrm{\Pi }}^2(\tau ,r)}{\mathrm{\Pi }^2(\tau ,r)}}+{\displaystyle \frac{\ddot{\mathrm{\Pi }}(\tau ,r)}{\mathrm{\Pi }(\tau ,r)}}{\displaystyle \frac{\mathrm{\Pi }^{\prime \prime }(\tau ,r)}{\mathrm{\Pi }(\tau ,r)}})\}.`$ (A.41)
## Appendix B Appendix to Chapter 4
### B.1 Details to section 4.2
In this appendix, the contraction of Lorentz and color indices appearing in the expressions for the Feynman diagrams Figs. 4.2, 4.4, 4.5 is performed. For diagrams containing TLM fluctuations and hence $`P_{\mu \nu }^T`$ as defined in Eq. (4.11), the following formulae will be useful:
$$\begin{array}{cc}\hfill P_{\mu \nu }^T\left(q\right)q^\nu & =0\hfill \\ \hfill P_{\mu \nu }^T\left(q\right)g^{\mu \nu }& =2\hfill \\ \hfill P_{\mu \nu }^T\left(q\right)p^\mu k^\nu & =𝒑𝒌\frac{\left(𝒒𝒑\right)\left(𝒒𝒌\right)}{𝒒^2}\hfill \\ \hfill P_{\mu \nu }^T\left(q\right)p^\mu p^\nu & =𝒑^2\frac{\left(𝒒𝒑\right)^2}{𝒒^2}.\hfill \end{array}$$
(B.1)
Concerning $`\mathrm{\Delta }P^{\mathrm{HH}}`$. The color indices $`a`$, $`b`$, $`c`$ and $`d`$ take the values 1,2 only because we chose the particles propagating in the loops to be TLH-modes. Index $`f`$ is summed over 1,2,3.
$`\left[\epsilon _{adf}\epsilon _{fbc}\left(g^{\mu \nu }g^{\rho \sigma }g^{\mu \rho }g^{\nu \sigma }\right)+\epsilon _{abf}\epsilon _{fdc}\left(g^{\mu \sigma }g^{\nu \rho }g^{\mu \rho }g^{\nu \sigma }\right)+\epsilon _{acf}\epsilon _{fdb}\left(g^{\mu \sigma }g^{\nu \rho }g^{\mu \nu }g^{\sigma \rho }\right)\right]`$
$`\left(\delta _{ab}\right)\left(g_{\mu \nu }{\displaystyle \frac{k_\mu k_\nu }{m^2}}\right)\left(\delta _{cd}\right)\left(g_{\rho \sigma }{\displaystyle \frac{p_\rho p_\sigma }{m^2}}\right)`$
$`=\epsilon _{acf}\epsilon _{acf}\left(g^{\mu \nu }g^{\rho \sigma }g^{\mu \rho }g^{\nu \sigma }g^{\mu \sigma }g^{\nu \rho }+g^{\mu \nu }g^{\rho \sigma }\right)\left(g_{\mu \nu }{\displaystyle \frac{k_\mu k_\nu }{m^2}}\right)\left(g_{\rho \sigma }{\displaystyle \frac{p_\rho p_\sigma }{m^2}}\right)`$ (B.2)
$`=4\left(g^{\mu \nu }g^{\rho \sigma }g^{\mu \rho }g^{\nu \sigma }\right)\left(g_{\mu \nu }{\displaystyle \frac{k_\mu k_\nu }{m^2}}\right)\left(g_{\rho \sigma }{\displaystyle \frac{p_\rho p_\sigma }{m^2}}\right)`$
$`=4\left(123{\displaystyle \frac{p^2}{m^2}}3{\displaystyle \frac{k^2}{m^2}}+{\displaystyle \frac{p^2k^2}{m^4}}{\displaystyle \frac{\left(pk\right)^2}{m^4}}\right)`$
Concerning $`\mathrm{\Delta }P^{\mathrm{MH}}`$. Since $`k`$ and $`p`$ are associated with a TLM-mode and a TLH-mode respectively, we have to set $`a=b=3`$ and sum over $`c,d=1,2`$. Eqs. (B.1) have to be used.
$`\left[\epsilon _{adf}\epsilon _{fbc}\left(g^{\mu \nu }g^{\rho \sigma }g^{\mu \rho }g^{\nu \sigma }\right)+\epsilon _{abf}\epsilon _{fdc}\left(g^{\mu \sigma }g^{\nu \rho }g^{\mu \rho }g^{\nu \sigma }\right)+\epsilon _{acf}\epsilon _{fdb}\left(g^{\mu \sigma }g^{\nu \rho }g^{\mu \nu }g^{\sigma \rho }\right)\right]`$
$`\left(\delta _{cd}\right)\left(g_{\rho \sigma }{\displaystyle \frac{p_\rho p_\sigma }{m^2}}\right)\left(\delta _{ab}\right)P_{\mu \nu }^T\left(k\right)`$
$`=\epsilon _{acf}\epsilon _{acf}\left(g^{\mu \nu }g^{\rho \sigma }g^{\mu \rho }g^{\nu \sigma }g^{\mu \sigma }g^{\nu \rho }+g^{\mu \nu }g^{\sigma \rho }\right)\left(g_{\rho \sigma }{\displaystyle \frac{p_\rho p_\sigma }{m^2}}\right)P_{\mu \nu }^T\left(k\right)`$
$`=4\left(g^{\mu \nu }g^{\rho \sigma }g^{\mu \rho }g^{\nu \sigma }\right)\left(g_{\rho \sigma }{\displaystyle \frac{p_\rho p_\sigma }{m^2}}\right)P_{\mu \nu }^T\left(k\right)`$ (B.3)
$`=4\left(\left(3{\displaystyle \frac{p^2}{m^2}}\right)g^{\mu \nu }+{\displaystyle \frac{p^\mu p^\nu }{m^2}}\right)P_{\mu \nu }^T\left(k\right)`$
$`=4\left(6+2{\displaystyle \frac{p^2}{m^2}}+{\displaystyle \frac{𝒑^2}{m^2}}{\displaystyle \frac{\left(𝒑𝒌\right)^2}{m^2𝒌^2}}\right)`$
Concerning $`\mathrm{\Delta }P^{\mathrm{MHH}}`$. The momentum $`q`$ is associated with a TLM fluctuation, and the TLH fluctuations carry momenta $`p`$ and $`k`$. Thus we have $`f,g=3`$ and $`a,b,c,d=1,2`$, and hence
$$\epsilon _{acf}\epsilon _{bdg}\delta _{ab}\delta _{cd}\delta _{fg}=\underset{a=1}{\overset{2}{}}\underset{c=1}{\overset{2}{}}\epsilon _{ac3}\epsilon _{ac3}=2.$$
(B.4)
The Lorentz structure resulting from the propagators and polarization tensors is
$`\left[g^{\mu \rho }\left(kp\right)^\lambda +g^{\rho \lambda }\left(qk\right)^\mu +g^{\lambda \mu }\left(pq\right)^\rho \right]\left[g^{\nu \sigma }\left(pk\right)^\kappa +g^{\sigma \kappa }\left(kq\right)^\nu +g^{\kappa \nu }\left(qp\right)^\sigma \right]`$
$`\left(g_{\mu \nu }{\displaystyle \frac{p_\mu p_\nu }{m^2}}\right)\left(g_{\rho \sigma }{\displaystyle \frac{k_\rho k_\sigma }{m^2}}\right)P_{\lambda \kappa }^T\left(q\right).`$ (B.5)
Contracting the Lorentz indices and inserting $`P_{\lambda \kappa }^T\left(q\right)`$ by using the rules (B.1) yields
$`2\left(𝒌^2{\displaystyle \frac{\left(𝒌𝒒\right)^2}{𝒒^2}}\right)2\left(𝒑^2{\displaystyle \frac{\left(𝒑𝒒\right)^2}{𝒒^2}}\right)+6\left(𝒌𝒑{\displaystyle \frac{\left(𝒌𝒒\right)\left(𝒑𝒒\right)}{𝒒^2}}\right)`$ (B.6)
$`+2k^24kq4pq+2p^2+4q^2`$
$`+{\displaystyle \frac{1}{m^2}}[(𝒑^2{\displaystyle \frac{\left(𝒑𝒒\right)^2}{𝒒^2}})(k^2+q^2)+(𝒌^2{\displaystyle \frac{\left(𝒌𝒒\right)^2}{𝒒^2}})(p^2+q^2)4(𝒌𝒑{\displaystyle \frac{\left(𝒌𝒒\right)\left(𝒑𝒒\right)}{𝒒^2}})kp`$
$`4\left(kp\right)^2+4\left(kp\right)\left(kq\right)2\left(kq\right)^2+4\left(kp\right)\left(pq\right)2\left(pq\right)^2]`$
$`+{\displaystyle \frac{1}{m^4}}\left[\left(𝒑^2{\displaystyle \frac{\left(𝒑𝒒\right)^2}{𝒒^2}}\right)\left(kq\right)^2\left(𝒌^2{\displaystyle \frac{\left(𝒌𝒒\right)^2}{𝒒^2}}\right)\left(pq\right)^2+2\left(𝒌𝒑{\displaystyle \frac{\left(𝒌𝒒\right)\left(𝒑𝒒\right)}{𝒒^2}}\right)\left(kq\right)\left(pq\right)\right].`$
Now we may use energy momentum conservation and set $`q=pk`$. Moreover, we notice that
$$𝒌^2\frac{\left(𝒌^2+𝒌𝒑\right)^2}{\left(𝒌+𝒑\right)^2}=𝒑^2\frac{\left(𝒑^2+𝒌𝒑\right)^2}{\left(𝒌+𝒑\right)^2}=\frac{\left(𝒌^2+𝒌𝒑\right)\left(𝒑^2+𝒌𝒑\right)}{\left(𝒌+𝒑\right)^2}𝒌𝒑=\frac{𝒌^2𝒑^2\left(𝒌𝒑\right)^2}{\left(𝒌+𝒑\right)^2}.$$
(B.7)
This leads to the following:
$`10{\displaystyle \frac{\left(𝒌𝒑\right)^2𝒌^2𝒑^2}{\left(𝒌+𝒑\right)^2}}+10k^2+16kp+10p^2`$ (B.8)
$`+{\displaystyle \frac{1}{m^2}}\left(2k^48k^2kp16\left(kp\right)^28kpp^22p^4{\displaystyle \frac{\left(𝒌𝒑\right)^2𝒌^2𝒑^2}{\left(𝒌+𝒑\right)^2}}\left(3k^2+8kp+3p^2\right)\right)`$
$`+{\displaystyle \frac{1}{m^4}}{\displaystyle \frac{\left(𝒌𝒑\right)^2𝒌^2𝒑^2}{\left(𝒌+𝒑\right)^2}}\left(k^4+4k^2kp+4\left(kp\right)^2+2k^2p^2+4kpp^2+p^4\right).`$
We are interested in only one contribution from this diagram, namely $`\mathrm{\Delta }P_{\mathrm{VTT}}^{\mathrm{MHH}}`$. Fortunately, its calculation includes the two on mass shell conditions $`p^2=m^2`$ and $`k^2=m^2`$. They reduce the former expression to
$$16\left(m^2\frac{\left(kp\right)^2}{m^2}\right)\frac{𝒌^2𝒑^2\left(𝒌𝒑\right)^2}{\left(𝒑+𝒌\right)^2}\left(8+4\frac{\left(kp\right)^2}{m^4}\right).$$
(B.9)
### B.2 Electric screening mass for TLM modes
There are two nonvanishing contributions to the polarization tensor $`\mathrm{\Pi }_{\mu \nu }\left(p\right)`$ for TLM modes on one-loop level. They are depicted in Figs. B.1 and B.2. We are especially interested in the component $`\mathrm{\Pi }_{00}`$ for vanishing external momentum, more precisely $`\mathrm{\Pi }_{00}\left(p_0=0,\left|𝒑\right|0\right)`$.
#### B.2.1 Tadpole contribution
According to the Feynman rules (see Sec. 4.1), the tadpole diagram Fig. B.1 corresponds to the expression
$`\mathrm{\Pi }_{\text{tadp.}}^{\mu \nu }\left(p\right)`$ $`={\displaystyle \frac{1}{i}}{\displaystyle \frac{d^4k}{\left(2\pi \right)^4}\left(\delta _{cd}\right)\left(g_{\rho \sigma }\frac{k_\rho k_\sigma }{m^2}\right)\left[\frac{i}{k^2m^2}+2\pi \delta \left(k^2m^2\right)n_B\left(\left|k_0\right|/T\right)\right]}`$
$`(ie^2)[\epsilon _{eab}\epsilon _{ecd}(g^{\mu \rho }g^{\nu \sigma }g^{\mu \sigma }g^{\nu \rho })+\epsilon _{eac}\epsilon _{edb}(g^{\mu \sigma }g^{\rho \nu }g^{\mu \nu }g^{\rho \sigma })`$ (B.10)
$`+\epsilon _{ead}\epsilon _{ebc}(g^{\mu \nu }g^{\sigma \rho }g^{\mu \rho }g^{\sigma \nu })]`$
$`=4e^2{\displaystyle \frac{d^4k}{\left(2\pi \right)^4}\left[g^{\mu \nu }\left(3\frac{k^2}{m^2}\right)+\frac{k^\mu k^\nu }{m^2}\right]\left[\frac{i}{k^2m^2}+2\pi \delta \left(k^2m^2\right)n_B\left(\left|k_0\right|/T\right)\right]}.`$
When integrating the loop momentum $`k`$, the following constraints have to be taken into account:
1. For both thermal and vacuum contribution the constraint on the center-of-mass energy in vertices,
$$\left|\left(p+k\right)^2\right|\left|\varphi \right|^2.$$
(B.11)
2. For the vacuum contribution additionally the constraint on the offshellness of quantum fluctuations,
$$\left|k^2m^2\right|\left|\varphi \right|^2.$$
(B.12)
As for the calculations of two-loop diagrams in Sec. 4, massive vacuum fluctuations for $`e>\frac{1}{2}`$ are forbidden by Eq. (B.12), i. e. the vacuum contribution vanishes for $`e>\frac{1}{2}`$. Moreover, for $`e>\frac{1}{2}`$ and external momentum $`p=0`$, the thermal part also vanishes because in this case
$$\left|\left(p+k\right)^2\right|=\left|k^2\right|=m^2=\left(2e\right)^2\left|\varphi \right|^2>\left|\varphi \right|^2,$$
(B.13)
and Eq. (B.11) cannot be satisfied. Thus we have no contribution from the tadpole for $`e>1/2`$.
#### B.2.2 Nonlocal contribution
The nonlocal contribution to the polarization tensor is depicted in Fig. B.2. In formula, it reads
$`\mathrm{\Pi }_{\text{nonl.}}^{\mu \nu }\left(p\right)`$ $`={\displaystyle \frac{1}{2i}}{\displaystyle \frac{d^4k}{\left(2\pi \right)^4}e^2\epsilon _{ace}\left[g^{\mu \rho }\left(pk\right)^\lambda +g^{\rho \lambda }\left(kp+k\right)^\mu +g^{\lambda \mu }\left(pk+p\right)^\rho \right]}`$
$`\epsilon _{dbf}\left[g^{\sigma \nu }\left(kp\right)^\kappa +g^{\nu \kappa }\left(p+pk\right)^\sigma +g^{\kappa \sigma }\left(p+k+k\right)^\nu \right]`$
$`\left(\delta _{cd}\right)\left(g_{\rho \sigma }{\displaystyle \frac{k_\rho k_\sigma }{m^2}}\right)\left[{\displaystyle \frac{i}{k^2m^2}}+2\pi \delta \left(k^2m^2\right)n_B\left(\left|k_0\right|/T\right)\right]`$
$`\left(\delta _{ef}\right)\left(g_{\lambda \kappa }{\displaystyle \frac{\left(pk\right)_\lambda \left(pk\right)_\kappa }{m^2}}\right)`$
$`\left[{\displaystyle \frac{i}{\left(pk\right)^2m^2}}+2\pi \delta \left(\left(pk\right)^2m^2\right)n_B\left(\left|p_0k_0\right|/T\right)\right]`$
$`=ie^2{\displaystyle }{\displaystyle \frac{d^4k}{\left(2\pi \right)^4}}\{(2k^22{\displaystyle \frac{k^4}{m^2}}2kp+4{\displaystyle \frac{k^2kp}{m^2}}4{\displaystyle \frac{\left(kp\right)^2}{m^2}}+5p^2+2{\displaystyle \frac{k^2p^2}{m^2}}{\displaystyle \frac{p^4}{m^2}})g^{\mu \nu }`$
$`+\left(10+2{\displaystyle \frac{k^2}{m^2}}2{\displaystyle \frac{kp}{m^2}}3{\displaystyle \frac{p^2}{m^2}}+{\displaystyle \frac{p^4}{m^4}}\right)k^\mu k^\nu +\left(23{\displaystyle \frac{k^2}{m^2}}+{\displaystyle \frac{\left(kp\right)^2}{m^4}}+{\displaystyle \frac{p^2}{m^2}}\right)p^\mu p^\nu `$
$`+(5{\displaystyle \frac{k^2}{m^2}}+4{\displaystyle \frac{kp}{m^2}}{\displaystyle \frac{p^2kp}{m^4}})(k^\nu p^\mu +k^\mu p^\nu )\}`$
$`\left[{\displaystyle \frac{i}{k^2m^2}}+2\pi \delta \left(k^2m^2\right)n_b\left(\left|k_0\right|/T\right)\right]`$
$`\left[{\displaystyle \frac{i}{\left(pk\right)^2m^2}}+2\pi \delta \left(\left(pk\right)^2m^2\right)n_B\left(\left|p_0k_0\right|/T\right)\right]`$ (B.14)
Again, all contributions containing a vacuum TLH fluctuation vanish for $`e>1/2`$. The only nonvanishing contribution is the case of both fluctuations being thermal; it reads
$`\mathrm{\Pi }_{\begin{array}{c}\text{nonl.}\\ \text{therm.}\end{array}}^{\mu \nu }\left(p\right)`$ $`=ie^2{\displaystyle }{\displaystyle \frac{d^4k}{\left(2\pi \right)^2}}\{(2kp4{\displaystyle \frac{\left(kp\right)^2}{m^2}}+7p^2{\displaystyle \frac{p^4}{m^2}})g^{\mu \nu }+(122{\displaystyle \frac{kp}{m^2}}3{\displaystyle \frac{p^2}{m^2}}+{\displaystyle \frac{p^4}{m^4}})k^\mu k^\nu `$
$`+(6+4{\displaystyle \frac{kp}{m^2}}{\displaystyle \frac{p^2kp}{m^4}})(k^\nu p^\mu +k^\mu p^\nu )+(5+{\displaystyle \frac{\left(kp\right)^2}{m^4}}+{\displaystyle \frac{p^2}{m^2}})p^\mu p^\nu \}`$
$`\delta \left(k^2m^2\right)n_B\left(\left|k_0\right|/T\right)\delta \left(\left(pk\right)^2m^2\right)n_B\left(\left|p_0k_0\right|/T\right).`$ (B.15)
We now set $`p_0=0`$ (keeping $`\left|𝒑\right|0`$) and get
$`\mathrm{\Pi }_{\begin{array}{c}\text{nonl.}\\ \text{therm.}\end{array}}^{\mu \nu }\left(p_0=0,𝒑\right)`$
$`=ie^2{\displaystyle }{\displaystyle \frac{d^4k}{\left(2\pi \right)^2}}\{(2𝒌𝒑4{\displaystyle \frac{\left(𝒌𝒑\right)^2}{m^2}}7𝒑^2{\displaystyle \frac{𝒑^4}{m^2}})g^{\mu \nu }+(12+2{\displaystyle \frac{𝒌𝒑}{m^2}}+3{\displaystyle \frac{𝒑^2}{m^2}}+{\displaystyle \frac{𝒑^4}{m^4}})k^\mu k^\nu `$
$`+(64{\displaystyle \frac{𝒌𝒑}{m^2}}{\displaystyle \frac{𝒑^2𝒌𝒑}{m^4}})(k^\nu p^\mu +k^\mu p^\nu )+(5+{\displaystyle \frac{\left(𝒌𝒑\right)^2}{m^4}}{\displaystyle \frac{𝒑^2}{m^2}})p^\mu p^\nu \}`$
$`\delta \left(k^2m^2\right)\delta \left(\left(pk\right)^2m^2\right)\left[n_B\left(\left|k_0\right|/T\right)\right]^2`$ (B.16)
The two $`\delta `$-functions can for $`p0`$ be rewritten as
$`\delta \left(k^2m^2\right)\delta \left(\left(pk\right)^2m^2\right)=\delta \left(k^2m^2\right)\delta \left(p^22pk\right)`$
$`={\displaystyle \frac{1}{2\sqrt{𝒌^2+m^2}}}[\delta (k_0\sqrt{𝒌^2+m^2})\delta (p_0^2𝒑^22p_0\sqrt{𝒌^2+m^2}+2𝒑𝒌)`$
$`+\delta (k_0+\sqrt{𝒌^2+m^2})\delta (p_0^2𝒑^2+2p_0\sqrt{𝒌^2+m^2}+2𝒑𝒌)]`$
$`={\displaystyle \frac{1}{2\sqrt{\left|𝒌\right|^2+m^2}}}[\delta (k_0\sqrt{\left|𝒌\right|^2+m^2})\delta (p_0^2|𝒑|^22p_0\sqrt{\left|𝒌\right|^2+m^2}+2\left|𝒑\right|\left|𝒌\right|\mathrm{cos}\theta )`$
$`+\delta (k_0+\sqrt{\left|𝒌\right|^2+m^2})\delta (p_0^2|𝒑|^2+2p_0\sqrt{\left|𝒌\right|^2+m^2}+2\left|𝒑\right|\left|𝒌\right|\mathrm{cos}\theta )].`$ (B.17)
In the case $`p_0=0`$ (with $`\left|𝒑\right|0`$) the former expression reduces to
$`\delta \left(k^2m^2\right)\delta \left(\left(pk\right)^2m^2\right)`$
$`={\displaystyle \frac{1}{2\sqrt{\left|𝒌\right|^2+m^2}}}[\delta (k_0\sqrt{\left|𝒌\right|^2+m^2})\delta \left(2\right|𝒑\left|\right|𝒌|\mathrm{cos}\theta |𝒑|^2)`$ (B.18)
$`+\delta (k_0+\sqrt{\left|𝒌\right|^2+m^2})\delta \left(2\right|𝒑\left|\right|𝒌|\mathrm{cos}\theta |𝒑|^2)]`$
$`={\displaystyle \frac{1}{4\left|𝒑\right|\left|𝒌\right|\sqrt{\left|𝒌\right|^2+m^2}}}\delta \left(\mathrm{cos}\theta {\displaystyle \frac{\left|𝒑\right|}{2\left|𝒌\right|}}\right)\left[\delta \left(k_0\sqrt{\left|𝒌\right|^2+m^2}\right)+\delta \left(k_0+\sqrt{\left|𝒌\right|^2+m^2}\right)\right].`$
Inserting Eq. (B.2.2) into Eq. (B.2.2) yields
$$\begin{array}{cc}& \mathrm{\Pi }_{\begin{array}{c}\text{nonl.}\\ \text{therm.}\end{array}}^{\mu \nu }\left(p_0=0,𝒑\right)\hfill \\ & =\frac{ie^2}{\left(2\pi \right)^2}_{\mathrm{}}^{\mathrm{}}𝑑k_0_{p/2}^{\mathrm{}}k^2𝑑k_{S_2}𝑑\mathrm{\Omega }\left[n_B\left(\sqrt{k^2+m^2}/T\right)\right]^2\hfill \\ & \{(8p^22\frac{p^4}{m^2})g^{\mu \nu }+(12+4\frac{p^2}{m^2}+\frac{p^4}{m^4})k^\mu k^\nu \hfill \\ & +(62\frac{p^2}{m^2}\frac{1}{2}\frac{p^4}{m^4})(k^\nu p^\mu +k^\mu p^\nu )+(5\frac{p^2}{m^2}+\frac{1}{4}\frac{p^4}{m^4})p^\mu p^\nu \}\hfill \\ & \frac{\delta \left(\mathrm{cos}\theta \frac{p}{2k}\right)}{4pk\sqrt{k^2+m^2}}\left[\delta \left(k_0\sqrt{k^2+m^2}\right)+\delta \left(k_0+\sqrt{k^2+m^2}\right)\right]\hfill \end{array}$$
(B.19)
Only the terms proportional to $`g^{\mu \nu }`$ or $`k^\mu k^\nu `$ contribute to $`\mathrm{\Pi }^{00}`$. The result is
$$\begin{array}{cc}& \mathrm{\Pi }_{\begin{array}{c}\text{nonl.}\\ \text{therm.}\end{array}}^{00}\left(p_0=0,𝒑\right)\hfill \\ & =\frac{ie^2}{2\pi }\left(4p+\frac{p^3}{m^2}\right)_{p/2}^{\mathrm{}}𝑑k\frac{k}{\sqrt{k^2+m^2}}n_B^2\left(\sqrt{k^2+m^2}/T\right)\hfill \\ & +\frac{ie^2}{4\pi }\left(\frac{12}{p}+4\frac{p}{m^2}+\frac{p^3}{m^4}\right)_{p/2}^{\mathrm{}}𝑑kk\sqrt{k^2+m^2}n_B^2\left(\sqrt{k^2+m^2}/T\right).\hfill \end{array}$$
(B.20)
In the limit $`\left|𝒑\right|0`$ this expression diverges,
$$\begin{array}{cc}\hfill \mathrm{\Pi }_{\begin{array}{c}\text{nonl.}\\ \text{therm.}\end{array}}^{00}\left(p_0=0,𝒑\right)& \underset{\left|𝒑\right|0}{\overset{}{}}\frac{3ie^2}{\pi }\frac{1}{\left|𝒑\right|}_0^{\mathrm{}}𝑑kk\sqrt{k^2+m^2}n_B^2\left(\sqrt{k^2+m^2}/T\right).\hfill \end{array}$$
(B.21) |
warning/0506/hep-ph0506048.html | ar5iv | text | # Semileptonic 𝐵→𝜋 Decays from an Omnès Improved Nonrelativistic Constituent Quark Model.
## I Introduction
Exclusive semileptonic decays of $`B`$-mesons are of great interest, since they can be used to determine the Cabibbo-Kobayashi-Maskawa (CKM) matrix elements $`|V_{ub}|`$ and $`|V_{cb}|`$. In the latter case, heavy quark symmetry greatly simplifies the theoretical understanding of the hadronic transition matrix elements and thus the overall theoretical uncertainty on the decay process is under control pdg . The measurement of the exclusive semileptonic decay $`B\pi l^+\nu _l`$ by the CLEO Collaboration Exp\_96 ; Exp\_03 can be used to determine the CKM parameter $`|V_{ub}|`$. This exclusive method provides an important alternative to the extraction of $`|V_{ub}|`$ from inclusive measurements of $`BX_ul^+\nu _l`$. For semileptonic decays of charmed or bottom mesons into light mesons there are no flavor symmetries to constrain the hadronic matrix elements, and as a result, the errors on $`|V_{ub}|`$ are currently dominated<sup>1</sup><sup>1</sup>1 The current best value for $`|V_{cd}|`$ comes from neutrino production of charm off valence $`d`$ quarks (with the cross section from perturbative QCD), rather than from semileptonic $`D`$ decays. by theoretical uncertainties pdg . An accurate determination of $`|V_{ub}|`$ with well–understood uncertainties remains one of the fundamental priorities for heavy flavor physics.
The transition amplitude for the exclusive semileptonic $`bu`$ decays factorizes into leptonic and hadronic parts. The hadronic matrix elements contain the non-perturbative, strong–interaction effects and have been extensively evaluated within different approaches. Thus, several lattice QCD (first in the quenched approximation, Latt\_96 ; Latt\_98 ; Latt\_98\_nrqcd ; Latt\_00 ; Latt\_01 ; Latt\_01\_bis , and more recently using dynamical configurations Latt\_04\_HPQCD ; Latt\_04\_FNLAB ), light-cone sum rule (LCSR) LCSR\_98 ; LCSR\_00 ; LCSR\_Hqet\_01 ; LCSR\_02 ; LCSR\_03 and constituent quark model QM\_85 ; QM\_89 ; QM\_90\_IW ; QM\_90 ; QM\_95 ; QM\_96 ; QM\_97 ; QM\_00 ; QM\_02 calculations have been carried out in recent years. Each of the above methods has only a limited range of applicability, namely: LCSR are suitable for describing the low squared momentum transfer ($`q^2`$) region of the form factors, while Lattice QCD, because of the limitation on the magnitude of spatial momentum components, provides results only for the high $`q^2`$ region. Constituent quark models may give the form factors in the full $`q^2`$ range, but they are not closely related to the QCD Lagrangian<sup>2</sup><sup>2</sup>2A rigorous derivation of this approach as an effective theory of QCD in the non–perturbative regime has not been obtained. and therefore have input parameters which are not directly measurable and might not be of fundamental significance. Thus, it is evident that a combination of various methods is required.
Watson’s theorem for the $`B\pi l^+\nu _l`$ process allows one to write a dispersion relation for each of the form factors entering in the hadronic matrix element. This procedure leads to the so-called Omnès representation Omnes , which can be used to constrain the $`q^2`$ dependence of the form factors from the elastic $`\pi B\pi B`$ scattering amplitudes Omnes\_01 . In Ref. Omnes\_01 , once-subtracted dispersion relations were used, and though promising results were found, they suffered from sizeable uncertainties because of imprecise knowledge of the $`\pi B\pi B`$ phase shifts far from threshold. A recent re-analysis of the Omnès representation in this context Omnes\_05 , has shown that the use of multiply-subtracted dispersion relations considerably diminishes the form factor dependence on the elastic $`\pi B\pi B`$ scattering amplitudes at high energies, and more importantly points out that the Omnès representation of the form factors can be used to combine predictions from various methods in different $`q^2`$ regions.
In this paper we study the semileptonic $`B\pi l^+\nu _l`$ decay. We take advantage of the findings of Ref. Omnes\_05 and use the predictions of LCSR calculations at $`q^2=0`$ to extend the predictions of a simple nonrelativistic constituent quark model (NRCQM) from its region of applicability (near $`q_{\mathrm{max}}^2=(m_Bm_\pi )^2`$) to all $`q^2`$ values accessible in the physical decay. We also use the available lattice QCD data to test our approach. We use a Monte Carlo procedure to find theoretical error bands for the form factors and the decay width. From our estimate of the decay width and the branching ratio measurement of Ref. Exp\_03 we obtain
$$|V_{ub}|_{\mathrm{this}\mathrm{work}}=0.0034\pm 0.0003(\mathrm{exp})\pm 0.0007(\mathrm{theory})$$
(1)
To further test this simple framework, we also study the $`D\pi `$ and $`DK`$ decays, for which there exist precise experimental data and for which the relevant CKM matrix elements ($`|V_{cd}|`$ and $`|V_{cs}|`$) are also well known. We find
$$f_\pi ^+(0)=0.63\pm 0.02,f_K^+(0)=0.79\pm 0.01,\frac{f_\pi ^+(0)}{f_K^+(0)}=0.80\pm 0.03,\frac{(D^0\pi ^{}e^+\nu _e)}{(D^0K^{}e^+\nu _e)}|_{\mathrm{this}\mathrm{work}}=0.079\pm 0.008$$
(2)
The plan of this paper is as follows. After this introduction, we study the semileptonic $`B\pi `$ decay in Sect. II. First we set up the form factor decomposition (Subsect. II.1), discuss the valence quark approximation (Subsect. II.2) and the role played by the $`B^{}`$ resonance (Subsect. II.3). The Omnès dispersion relation and its application to this decay is addressed in Subsect. II.4 and in the Appendix. Finally, in Subsect. II.5 we use our framework to determine $`|V_{ub}|`$, paying special attention to estimating the uncertainties of the determination. In Sect. III we study the $`D\pi `$ and $`DK`$ semileptonic decays and finally in Sect. IV we present our conclusions.
## II Semileptonic $`B\pi `$ Decays
### II.1 Differential Decay Width and Form Factor Decomposition
Using Lorentz, parity, and time-reversal invariance, the matrix element for the semileptonic $`B^0\pi ^{}l^+\nu _l`$ decay can be parametrized in terms of two invariant and dimensionless form factors as<sup>3</sup><sup>3</sup>3Note that the axial current does not contribute to transitions between pseudoscalar mesons.
$$\pi (p_\pi )|V^\mu |B(p_B)=\left(p_B+p_\pi q\frac{m_B^2m_\pi ^2}{q^2}\right)^\mu f^+(q^2)+q^\mu \frac{m_B^2m_\pi ^2}{q^2}f^0(q^2)$$
(3)
where $`q^\mu =(p_Bp_\pi )^\mu `$ is the four momentum transfer and $`m_B=5279.4`$ MeV and $`m_\pi =139.57`$ MeV are the $`B^0`$ and $`\pi ^{}`$ masses, respectively. The physical meaning of the form factors is clear in the helicity basis, in which $`f^+`$ ($`f^0`$) corresponds to a transition amplitude with $`1^{}`$ ($`0^+`$) spin–parity quantum numbers in the center of mass of the lepton pair. For massless leptons ($`l=e`$ or $`\mu `$), the total decay rate is given by
$$\mathrm{\Gamma }\left(B^0\pi ^{}l^+\nu _l\right)=\frac{G_F^2|V_{ub}|^2}{192\pi ^3m_B^3}_0^{q_{\mathrm{max}}^2}𝑑q^2\left[\lambda (q^2)\right]^{\frac{3}{2}}|f^+(q^2)|^2$$
(4)
with $`q_{\mathrm{max}}^2=(m_Bm_\pi )^2`$, $`G_F=1.16637\times 10^5\mathrm{GeV}^2`$ and $`\lambda (q^2)=(m_B^2+m_\pi ^2q^2)^24m_B^2m_\pi ^2=4m_B^2|\stackrel{}{p}_\pi |^2`$, with $`\stackrel{}{p}_\pi `$ the pion three-momentum in the $`B`$ rest frame.
Measurements of the $`B^0`$ lifetime, $`\tau _{B^0}=(1.536\pm 0.014)\times 10^{12}\mathrm{s}`$ and of the $`B^0\pi ^{}l^+\nu _l`$ branching fraction, $`_{\mathrm{exp}}(B^0\pi ^{}l^+\nu _l)=(1.33\pm 0.22)\times 10^4`$ pdg lead to
$$\mathrm{\Gamma }_{\mathrm{exp}}\left(B^0\pi ^{}l^+\nu _l\right)=(8.7\pm 1.5)\times 10^7\mathrm{s}^1=(5.7\pm 1.0)\times 10^{14}\mathrm{MeV},l=e\mathrm{or}\mu $$
(5)
### II.2 Nonrelativistic Constituent Quark Model: Valence Quark Contribution
Within the spectator approximation, considering only the valence quark contribution and assuming that the $`B`$ and $`\pi `$ mesons are $`S`$-wave quark-antiquark bound states, a NRCQM (with constituent quark masses $`m_b`$ and $`m_l=m_u=m_d`$) predicts NRCQM\_def :
$`{\displaystyle \frac{\pi (E_\pi ,\stackrel{}{q})\left|V^\mu \right|B(m_B,\stackrel{}{0})^{\mathrm{val}}}{\sqrt{4m_BE_\pi }}}`$ $`=`$ $`{\displaystyle \frac{d^3l}{4\pi }\sqrt{\frac{E_b(\stackrel{}{l})+m_b}{2E_b(\stackrel{}{l})}}\sqrt{\frac{E_u(\stackrel{}{l}+\stackrel{}{q})+m_u}{2E_u(\stackrel{}{l}+\stackrel{}{q})}}\varphi _{\mathrm{rel}}^B(|\stackrel{}{l}|)\varphi _{\mathrm{rel}}^\pi (|\stackrel{}{l}+\frac{m_{sp}}{m_u+m_{sp}}\stackrel{}{q}|)𝒱^\mu (\stackrel{}{l},\stackrel{}{q})}`$
$`𝒱^\mu (\stackrel{}{l},\stackrel{}{q})`$ $`=`$ $`\left(\begin{array}{c}1+\frac{\stackrel{}{l}^{\mathrm{\hspace{0.17em}\hspace{0.17em}2}}+\stackrel{}{l}\stackrel{}{q}}{(E_b(\stackrel{}{l})+m_b)(E_u(\stackrel{}{l}+\stackrel{}{q})+m_u)}\\ \\ \frac{\stackrel{}{l}}{E_b(\stackrel{}{l})+m_b}\frac{\stackrel{}{l}+\stackrel{}{q}}{E_u(\stackrel{}{l}+\stackrel{}{q})+m_u)}\end{array}\right)`$ (9)
with $`E_\pi =\sqrt{m_\pi ^2+\stackrel{}{q}^{\mathrm{\hspace{0.17em}\hspace{0.17em}2}}}`$, $`E_{b,u}(\stackrel{}{k})=\sqrt{m_{b,u}^2+\stackrel{}{k}^{\mathrm{\hspace{0.17em}2}}}`$, $`m_{sp}`$ the spectator quark mass ($`m_d`$ in this case) and $`\varphi _{\mathrm{rel}}^{B,\pi }(k)`$ the Fourier transforms of the radial coordinate space $`B,\pi `$ meson wave functions, which describe the relative dynamics of the quark-antiquark pair<sup>4</sup><sup>4</sup>4They are normalized to $`_0^+\mathrm{}𝑑kk^2|\varphi _{\mathrm{rel}}^{B,\pi }(k)|^2=1`$.
To evaluate the coordinate space wave function we have used several nonrelativistic quark-antiquark interactions. Their general structure is as follows QM\_pot\_1 ; QM\_pot\_2
$$V_{ij}^{q\overline{q}}(r)=\frac{\kappa \left(1e^{r/r_c}\right)}{r}+\lambda r^p\mathrm{\Lambda }+\left\{a_0\frac{\kappa }{m_im_j}\frac{e^{r/r_0}}{rr_0^2}+\frac{2\pi }{3m_im_j}\kappa ^{}\left(1e^{r/r_c}\right)\frac{e^{r^2/x_0^2}}{\pi ^{\frac{3}{2}}x_0^3}\right\}\stackrel{}{\sigma }_i\stackrel{}{\sigma }_j$$
(10)
with $`\stackrel{}{\sigma }`$ the spin Pauli matrices, $`m_i`$ the constituent quark masses and
$$x_0(m_i,m_j)=A\left(\frac{2m_im_j}{m_i+m_j}\right)^B$$
(11)
The potentials considered differ in the form factors used for the hyperfine terms, the power of the confining term ($`p=1`$, as suggested by lattice QCD calculations GM84 , or $`p=2/3`$ which gives the correct asymptotic Regge trajectories for mesons Fabre88 ), or the use of a form factor in the one gluon exchange Coulomb potential. All interactions have been adjusted to reproduce the light ($`\pi `$, $`\rho `$, $`K`$, $`K^{}`$, etc.) and heavy-light ($`D`$, $`D^{}`$, $`B`$, $`B^{}`$, etc.) meson spectra and lead to precise predictions for the charmed and bottom baryon ($`\mathrm{\Lambda }_{c,b}`$, $`\mathrm{\Sigma }_{c,b}`$, $`\mathrm{\Sigma }_{c,b}^{}`$, $`\mathrm{\Xi }_{c,b}`$, $`\mathrm{\Xi }_{c,b}^{}`$, $`\mathrm{\Xi }_{c,b}^{}`$, $`\mathrm{\Omega }_{c,b}`$ and $`\mathrm{\Omega }_{c,b}^{}`$) masses QM\_pot\_2 ; QM\_hqs and for the semileptonic $`\mathrm{\Lambda }_b^0\mathrm{\Lambda }_c^+l^{}\overline{\nu }_l`$ and $`\mathrm{\Xi }_b^0\mathrm{\Xi }_c^+l^{}\overline{\nu }_l`$ QM\_slp decays.
Typical NRCQM valence quark predictions for the $`f^+`$ and $`f^0`$ form factors are depicted in Fig. 1. The AL1 potential from Ref. QM\_pot\_2 has been used<sup>5</sup><sup>5</sup>5The sensitivity of the results to the quark-antiquark nonrelativistic interaction will be discussed in detail later. and for comparison quenched lattice results are also plotted. Preliminary unquenched lattice calculations have been presented recently Latt\_04\_HPQCD ; Latt\_04\_FNLAB , but no significant difference between quenched and unquenched calculations is observed Has04 , within relatively large statistical errors. In addition, LCSR provide accurate and theoretically well founded results in the $`q^2=0`$ region. Thus, we have a LCSR value LCSR\_00
$$\mathrm{LCSR}:f^+(0)=0.28\pm 0.05$$
(12)
which is also plotted in Fig. 1.
Fig. 1 clearly shows the deficiencies of the NRCQM valence quark description of the $`B\pi l^+\nu _l`$ semileptonic decay. It fails over the full range of $`q^2`$ values. Close to $`q_{\mathrm{max}}^2`$, where the nonrelativistic Schrödinger equation should work best, the influence of the $`B^{}`$ resonance is clearly visible QM\_90\_IW . At the opposite end, close to $`q^2=0`$, where $`|\stackrel{}{q}|2.5\mathrm{GeV}`$, predictions from a nonrelativistic scheme are clearly not trustworthy. As a result, a value for the width $`\mathrm{\Gamma }_{\mathrm{NRCQM}}^{\mathrm{val}}\left(B^0\pi ^{}l^+\nu _l\right)=2.4\left(\frac{|V_{ub}|}{0.0032}\right)^2\times 10^{14}\mathrm{MeV}`$ is obtained, which is around a factor of two smaller than the CLEO measurement quoted in Eq. (5).
### II.3 Nonrelativistic Constituent Quark Model: $`B^{}`$ Resonance Contribution
A NRCQM description of the decay process should be feasible in the neighborhood of $`q_{\mathrm{max}}^2`$. Indeed, this is the case for the semileptonic $`BDl\overline{\nu }_l`$ and $`BD^{}l\overline{\nu }_l`$ decays, recently studied in Ref. NRCQM\_def with the same NRCQM as here. The difference here is that, as first pointed out in Ref. QM\_90\_IW , in the chiral limit and as $`m_b\mathrm{}`$, the decay $`B^0\pi ^{}l^+\nu _l`$ should be dominated near zero pion recoil by the effects of the $`B^{}`$ resonance, which is quite close to $`q_{\mathrm{max}}^2`$. In the picture of Ref. QM\_90\_IW , the one we will adopt here, the $`B^{}`$ contribution plays a role only near $`q_{\mathrm{max}}^2`$, since it is strongly suppressed by a soft hadronic vertex. This is in sharp contrast to phenomenological parameterizations of $`f^+`$ which assume it dominates over the full range accessible in the physical decay Latt\_96 . The $`B^{}`$ effects of the type considered here are not dual to the valence quark model form factors and must be added as a distinct coherent contribution to heavy quark decay near $`q_{\mathrm{max}}^2`$ QM\_90\_IW . We will focus on the $`f^+`$ form factor, which determines the decay width for massless leptons, and we evaluate the contribution to it from the diagram depicted in Fig. 2. It leads to a hadronic amplitude (normalizations as in Ref. NRCQM\_def )
$$iT^\mu =i\widehat{g}_{B^{}B\pi }(q^2)p_\pi ^\nu \left(i\frac{g_\nu ^\mu +q^\mu q_\nu /m_B^{}^2}{q^2m_B^{}^2}\right)i\sqrt{q^2}\widehat{f}_B^{}(q^2)$$
(13)
with $`m_B^{}=5325`$ MeV, and $`\widehat{f}_B^{}`$ and $`\widehat{g}_{B^{}B\pi }`$ the $`B^{}`$ decay constant and the strong $`B^{}B\pi `$ dimensionless coupling constant for a virtual $`B^{}`$ meson, respectively. On the $`B^{}`$ mass shell, the hadron matrix elements $`\widehat{f}_B^{}(q^2=m_B^{}^2)f_B^{}`$ and $`\widehat{g}_{B^{}B\pi }(q^2=m_B^{}^2)g_{B^{}B\pi }`$ reduce to the ordinary $`B^{}`$ decay constant and coupling of a pion to $`B`$ and $`B^{}`$ mesons. The latter is related, in the heavy quark limit, to $`\widehat{g}`$, the coupling of the vector and pseudoscalar heavy-light mesons to the pion defg ; hmchpt <sup>6</sup><sup>6</sup>6We use the normalization $`f_\pi 131\mathrm{MeV}`$.
$$g_{B^{}B\pi }=\left(\frac{2\widehat{g}\sqrt{m_Bm_B^{}}}{f_\pi }\right)(1+𝒪(1/m_b))$$
(14)
From Eq. (13) we get
$$f_{\mathrm{pole}}^+(q^2)=\frac{1}{2}\widehat{g}_{B^{}B\pi }(q^2)\frac{\sqrt{q^2}\widehat{f}_B^{}(q^2)}{m_B^{}^2q^2}$$
(15)
There is no direct experimental determination of $`g_{B^{}B\pi }`$, because there is no phase space for the decay $`B^{}B\pi `$. The available experimental results for $`D^{}D\pi `$ pdg can be related to $`g_{B^{}B\pi }`$, through heavy quark symmetry. There is no direct measurements of $`f_B^{}`$ either. In Ref. NRCQM\_def we computed, within the same NRCQM approach as the one outlined here, both $`g_{B^{}B\pi }`$ and $`f_B^{}`$, and we found a value of $`9.1\pm 0.9\mathrm{GeV}`$ for the product of both quantities, which appears in $`f_{\mathrm{pole}}^+`$ at $`q^2=m_B^{}^2`$. Lattice QCD simulations have measured $`f_B^{}`$ Latt\_fbstar and $`g_{B^{}B\pi }`$ g\_Latt\_03 to be
$`f_B^{}=190\pm 30\mathrm{MeV}g_{B^{}B\pi }=47\pm 5\pm 8\left[g_{B^{}B\pi }f_B^{}\right]_{\mathrm{Latt}\mathrm{QCD}}=8.9\pm 2.2\mathrm{GeV}`$ (16)
where we have added errors in quadrature. Thus the lattice prediction for the product $`g_{B^{}B\pi }f_B^{}`$ is in remarkable agreement, within $`3\%`$, with our NRCQM estimate in NRCQM\_def . In what follows we will use the value and error for the product estimated from the lattice data and use the NRCQM of Ref. NRCQM\_def to determine the $`q^2`$ dependence of $`\widehat{g}_{B^{}B\pi }(q^2)`$ and $`\widehat{f}_B^{}(q^2)`$, as we will discuss below. There are other recent estimates for $`g_{B^{}B\pi }`$ (g\_LCSR\_95 ; g\_QM\_99 ) and $`f_B^{}`$ (Be02 ), but given the existing uncertainties, all of them are compatible with the lattice values quoted in Eq. (16).
As mentioned above, we use the NRCQM framework to estimate the $`q^2`$ dependence of the product of $`\widehat{f}_B^{}(q^2)\widehat{g}_{B^{}B\pi }(q^2)`$. Since the NRCQM always uses on-shell meson wave functions, all $`q^2`$ dependence will arise from the kinematical factors relating the quark model matrix elements and the hadron form factors. For instance, from Eq. (13) of Ref. NRCQM\_def we find a rather mild $`q^2`$ dependence
$$\widehat{f}_B^{}(q^2)\sqrt[4]{q^2}=f_B^{}\sqrt{m_B^{}}$$
(17)
In the same manner, we use Eqs. (50) and (51) of Ref. NRCQM\_def to determine the $`q^2`$ dependence of $`\widehat{g}_{B^{}B\pi }`$, setting the $`B`$-meson four momentum $`P^\mu =(m_B,\stackrel{}{P}^{}=0)`$ and the $`B^{}`$-meson four momentum $`P^\mu =q^\mu =(m_BE_\pi ,\stackrel{}{P}=\stackrel{}{q})`$, and off–shell mass given by $`\sqrt{q^2}`$, $`|\stackrel{}{q}|=\sqrt{E_\pi ^2m_\pi ^2}`$, and $`E_\pi `$ determined from $`q^2`$ as usual ($`E_\pi =(m_B^2+m_\pi ^2q^2)/2m_B`$). Thus, finally we evaluate
$$f_{\mathrm{pole}}^+(q^2)=\frac{1}{2}G_B^{}(q^2)\frac{\sqrt[4]{q^2}}{\sqrt{m_B^{}}}\frac{m_B^{}\left[g_{B^{}B\pi }f_B^{}\right]_{\mathrm{Latt}\mathrm{QCD}}}{m_B^{}^2q^2}$$
(18)
where $`G_B^{}(q^2)=\widehat{g}_{B^{}B\pi }(q^2)/g_{B^{}B\pi }`$ is a dimensionless hadronic factor normalized to one at $`q^2=m_B^{}^2`$, which accounts for the $`q^2`$ dependence of $`BB^{}\pi `$ amplitude. In Fig. 3, we show the influence of the $`B^{}`$ resonance within our NRCQM and compare our results to those obtained by Isgur and Wise from the gaussian constituent quark model of Refs. QM\_89 ; QM\_90\_IW . Our model for the $`B^{}`$ contribution compares well to that of Ref. QM\_90\_IW , though the latter decreases faster owing to the use of a harmonic oscillator basis. The inclusion of $`f_{\mathrm{pole}}^+`$ clearly improves the simple valence quark contribution and leads to a reasonable description of the lattice data from $`q_{\mathrm{max}}^2`$ down to $`q^2`$ values around $`15\mathrm{GeV}^2`$. The low $`q^2`$ region is still poorly described within the current model since relativistic corrections there should be large.
The hadronic amplitude of Eq. (13) also leads to a small contribution to the $`f^0`$ form factor. Though it also improves the description for the highest $`q^2`$ values, it is not large enough and it is necessary to consider the influence of the lightest $`0^+`$ $`B`$-resonances Omnes\_01 (for instance a resonance around $`5660\mathrm{MeV}`$ Ci99 ).
### II.4 Omnès Representation
Here we use the Omnès representation of the $`f^+`$ form factor to combine the NRCQM predictions at high $`q^2`$ values, say above $`18\mathrm{GeV}^2`$, with the LCSR result at $`q^2=0`$. In this way we obtain the full $`q^2`$ dependence of the form factor and thus can determine the $`|V_{ub}|`$ CKM matrix element from the integrated semileptonic width. As shown in the Appendix, the $`(n+1)`$-subtracted Omnès representation for $`f^+`$ reads:
$`f^+(q^2)`$ $`=`$ $`\left({\displaystyle \underset{j=0}{\overset{n}{}}}\left[f^+(q_j^2)\right]^{\alpha _j(q^2)}\right)\mathrm{exp}\left\{_\delta (q^2;q_0^2,q_1^2,\mathrm{},q_n^2){\displaystyle \underset{k=0}{\overset{n}{}}}(q^2q_k^2)\right\},`$
$`_\delta (q^2;q_0^2,\mathrm{},q_n^2)`$ $`=`$ $`{\displaystyle \frac{1}{\pi }}{\displaystyle _{s_{\mathrm{th}}}^+\mathrm{}}{\displaystyle \frac{ds}{(sq_0^2)\mathrm{}(sq_n^2)}}{\displaystyle \frac{\delta (s)}{sq^2}},`$
$`\alpha _j(q^2)`$ $`=`$ $`{\displaystyle \underset{jk=0}{\overset{n}{}}}{\displaystyle \frac{q^2q_k^2}{q_j^2q_k^2}},`$ (19)
with $`q^2<s_{\mathrm{th}}=(m_B+m_\pi )^2`$ and $`q_0^2,\mathrm{},q_n^2]\mathrm{},s_{\mathrm{th}}[`$. This representation requires as an input the elastic $`\pi B\pi B`$ phase shift $`\delta (s)`$ in the $`J^P=1^{}`$ and isospin $`I=1/2`$ channel plus the form-factor at $`(n+1)`$ $`q^2`$-values ($`q_0^2,q_1^2,\mathrm{},q_n^2`$) below the $`\pi B`$ threshold.
We would like to stress that from a theoretical point of view the Omnès representation is derived from first principles: the well-established Mandelstam hypothesis Ma58 of maximum analyticity and Watson’s theorem Wa55
$`{\displaystyle \frac{f^+(s+\mathrm{i}ϵ)}{f^+(s\mathrm{i}ϵ)}}={\displaystyle \frac{T(s+\mathrm{i}ϵ)}{T(s\mathrm{i}ϵ)}}=e^{2\mathrm{i}\delta (s)},s>s_{\mathrm{th}},T(s)={\displaystyle \frac{8\pi \mathrm{i}s}{\lambda ^{\frac{1}{2}}(s)}}\left(e^{2i\delta (s)}1\right)`$ (20)
Omnès ideas have been used successfully to account for final state interactions in kaon decays Omnes\_01\_kaon <sup>7</sup><sup>7</sup>7There, however, multiple derivatives evaluated at a single point are used as input instead of subtractions for different $`q^2`$ values. and in Ref. Omnes\_01 , a once-subtracted Omnès representation (subtraction point $`q_0^2=0`$) was applied to the study of semi-leptonic $`B\pi `$ decays. In the latter work phase shifts were evaluated by solving the Bethe-Salpeter equation in the so-called *on-shell scheme* EJ99 , with a kernel determined by the direct tree level amplitude from the lowest order heavy meson chiral perturbation theory lagrangian hmchpt , together with the tree diagrams for $`B^{}`$ exchange which involve the leading interaction with coupling $`\widehat{g}`$. Such a model accommodates the $`B^{}`$ as a $`\pi B`$ bound state and should acceptably describe phase shifts close to threshold. It led to promising results for $`f^+`$ Omnes\_01 , but theoretical uncertainties on the form factor were not negligible, since to compute the Omnès factor $`_\delta `$ (Eq. (19)) requires elastic phase-shifts far from threshold<sup>8</sup><sup>8</sup>8Higher resonance effects on phase shifts cannot be neglected far from threshold. In particular the LCSR result at $`q^2=0`$ hints that at least an extra $`J^P=1^{}`$ resonance, located around $`6\mathrm{GeV}`$, has to be included in the once-subtracted Omnès relation scheme Omnes\_01 .. To include the effects of higher resonances on $`\delta (s)`$ requires input of the masses and couplings of such resonances. We therefore make many subtractions in the Omnès dispersion relation to suppress the impact of $`\delta (s)`$ at large $`s`$. This will leave a systematic effect in our results, but this should be less than that coming from the NRCQM plus $`B^{}`$ pole used as our main input.
As the number of subtractions increases the integration region relevant in Eq. (19) gets reduced and, if this number is large enough, only the phase shifts at or near threshold will be needed. Note that close to threshold the $`p`$-wave phase shift behaves as
$$\delta (s)=n_b\pi p^3a+\mathrm{}$$
(21)
where $`n_b`$ is the number of bound states in the channel (Levinson’s theorem MS70 ), $`p`$ is the $`\pi B`$ center of mass momentum and $`a`$ the corresponding scattering volume. In our case $`n_b=1`$ if we consider the $`B^{}`$ as a $`\pi B`$ bound state. Here, we will perform a large number of subtractions so that approximating $`\delta (s)\pi `$ in Eq. (19) will be justified. The Omnès factor $`_\delta `$ can then be evaluated analytically and we find for $`q^2<s_{\mathrm{th}}`$
$`f^+(q^2)`$ $``$ $`{\displaystyle \frac{1}{s_{\mathrm{th}}q^2}}{\displaystyle \underset{j=0}{\overset{n}{}}}\left[f^+(q_j^2)(s_{\mathrm{th}}q_j^2)\right]^{\alpha _j(q^2)},n1`$ (22)
Next we use the above formula to combine the LCSR result at $`q^2=0`$ and those obtained from our NRCQM in the high $`q^2`$ region and presented in the previous section. Thus we have used the $`f^+`$ NRCQM (valence $`+`$ pole) predictions for five $`q^2`$ values ranging from $`q_{\mathrm{max}}^2`$ down<sup>9</sup><sup>9</sup>9From Eq. (9), we see that the arguments of the meson wave function are $`|\stackrel{}{l}|`$ and $`|\stackrel{}{l}+\stackrel{}{q}/2|`$. For $`q^2=18\mathrm{GeV}^2`$, half of the transferred momentum, $`|\stackrel{}{q}|/2`$, is about $`0.4`$$`0.5\mathrm{GeV}`$, which is of the same order as $`\stackrel{}{l}^{\mathrm{\hspace{0.17em}2}}_{\varphi ^{B,\pi }}^{\frac{1}{2}}`$. Since the non-relativistic quark-antiquark interactions, $`V(r)`$, have been adjusted to reproduce the meson binding energies, they effectively incorporate some relativistic corrections and hence one might expect this effective nonrelativistic framework to provide reasonable meson wave functions for momenta of order $`\stackrel{}{l}^{\mathrm{\hspace{0.17em}2}}_{\varphi ^{B,\pi }}^{\frac{1}{2}}`$. This could explain why the NRCQM describes the lattice data (right panel of Fig. 3) from high values of $`q^2`$ down to values of $`q^2`$ even smaller than $`18\mathrm{GeV}^2`$. Nevertheless, we find it surprising that the nonrelativistic constituent quark model works as well as it does QM\_95 . to about $`18\mathrm{GeV}^2`$:
$$(q^2/\mathrm{GeV}^2,f^+(q^2))=\{\begin{array}{c}(23.574,4.1373),\hfill \\ (21.804,2.5821),\hfill \\ (21.116,2.1969),\hfill \\ (20.173,1.7916),\hfill \\ (18.290,1.2591)\hfill \end{array}$$
(23)
together with the LCSR result of Eq. (12) at $`q^2=0`$. When one uses a large number of subtractions, as is the case here, the $`\alpha _j`$ exponents become large and there are huge cancellations (note the normalization condition given in Eq. (52)). This is the reason why above, and to ensure numerical stability, we have quoted five significant digits for the NRCQM input. We are aware that uncertainties are larger than a precision of five digits, and we will carefully take this fact into account below. Results are shown in Fig. 4. As can be seen there, we obtain a simultaneous description of both lattice data in the high $`q^2`$ region and the LCSR prediction at $`q^2=0`$. In this way, starting from a nonrelativistic valence quark picture of the semileptonic process (Subsect. II.2) with all its obvious limitations, we have ended up with a realistic description of the relevant form factor for all $`q^2`$ values accessible in the physical decay.
A final remark concerns the use of the simplified Omnès representation of Eq. (22) instead of the exact one of Eq. (19). For instance, if we use five subtractions (we drop the NRCQM point at $`q^2=21.1\mathrm{GeV}^2`$) and the full Omnès representation<sup>10</sup><sup>10</sup>10We use the model of Ref. Omnes\_01 to obtain the phase-shifts. of Eq. (19), we find tiny differences from the results shown in Fig. 4. These differences are negligible (below $`1\%`$) above $`10\mathrm{GeV}^2`$, and though larger, still quite small (around $`5`$-$`7\%`$ at most in the $`5\mathrm{GeV}^2`$ region) below $`10\mathrm{GeV}^2`$.
### II.5 Determination of $`|V_{ub}|`$: Error Analysis
The CKM element $`|V_{ub}|`$ can be determined by comparing the experimental decay width (Eq. (5)) with the result of performing the phase space integration of Eq. (4) using the form factor $`f^+`$ determined in the previous subsection. Here, we will pay special attention to estimating the theoretical uncertainties. We have two main sources of theoretical errors:
1. *Uncertainties in the constituent quark-antiquark nonrelativistic interaction:* To estimate those, we will evaluate the spread of integrated widths obtained when five different potentials (AL1, AL2, AP1, AP2 and BD, in the notation of Ref. QM\_pot\_2 ) are considered. The forms and main characteristics of those potentials were discussed in Eqs. (10) and (11). As mentioned in Subsect. II.2, all interactions have been adjusted to reproduce the light and heavy–light meson spectra and lead to precise predictions for the charmed and bottom baryon masses QM\_pot\_2 and for the semileptonic $`\mathrm{\Lambda }_b^0\mathrm{\Lambda }_c^+l^{}\overline{\nu }_l`$ and $`\mathrm{\Xi }_b^0\mathrm{\Xi }_c^+l^{}\overline{\nu }_l`$ QM\_slp decays.
2. *Uncertainties on $`\left[g_{B^{}B\pi }f_B^{}\right]`$ and on the input to the multiply-subtracted Omnès representation:* Errors on $`\left[g_{B^{}B\pi }f_B^{}\right]`$, quoted in Eq. (16), affect the $`B^{}`$ pole contribution to $`f^+`$ (see Eq. (18)) and also induce uncertainties in the NRCQM prediction for the five points used as input to the Omnès representation in Eq. (22). Quark-antiquark potential uncertainties, discussed in the previous item, also induce uncertainties in the Omnès input. The errors on the $`q^2=0`$ data point (LCSR), quoted in Eq. (12), should also be taken into account.
To take these uncertainties into account, we proceed in two steps:
1. We fix the quark-antiquark potential to the AL1 interaction as in all previous subsections. By means of a Monte Carlo simulation, we generate a total of $`1000`$ $`(\left[g_{B^{}B\pi }f_B^{}\right]_{\mathrm{Latt}\mathrm{QCD}},f^+(0)_{\mathrm{LCSR}})`$ pairs<sup>11</sup><sup>11</sup>11We have checked that the errors quoted in the following are already stable when $`500`$ event simulations are performed. from an uncorrelated two dimensional gaussian distribution, with central values and standard deviations taken from Eqs. (16) and (12), respectively. For each of the $`1000`$ pairs we build up the six points that we use in our Omnès scheme (Eq. (22)) and thus find $`1000`$ different determinations of $`f^+`$ over the whole $`q^2`$ range accessible in the $`B\pi `$ decay. For each value of $`q^2`$, we discard the highest and lowest $`16\%`$ of the values obtained for the form factor, to leave a $`68\%`$ confidence level band which forms part of the theoretical uncertainty shown in Fig. 4. Since the output distributions are not gaussian in general, this accounts for possible skewness.
Performing the phase space integration for each of the $`1000`$ form factor samples and again discarding the highest and lowest $`16\%`$ of the values, we find
$$\frac{\mathrm{\Gamma }\left(B^0\pi ^{}l^+\nu _l\right)}{|V_{ub}|^2}=\left(0.50_{0.10}^{+0.14}\right)\times 10^8\mathrm{MeV}$$
(24)
2. We fix the $`(\left[g_{B^{}B\pi }f_B^{}\right]_{\mathrm{Latt}\mathrm{QCD}},f^+(0)_{\mathrm{LCSR}})`$ pair to their central values and compute the decay width with each of the five quark-antiquark interactions discussed above. From the spread of output values, we find
$$\frac{\mathrm{\Gamma }\left(B^0\pi ^{}l^+\nu _l\right)}{|V_{ub}|^2}=\left(0.50\pm 0.15\right)\times 10^8\mathrm{MeV}$$
(25)
Adding both sources of error in quadrature, we get
$$\frac{\mathrm{\Gamma }\left(B^0\pi ^{}l^+\nu _l\right)}{|V_{ub}|^2}=\left(0.50\pm 0.20\right)\times 10^8\mathrm{MeV}$$
(26)
and comparing to the measurement of the width in Eq. (5) we find
$$|V_{ub}|_{\mathrm{this}\mathrm{work}}=0.0034\pm 0.0003(\mathrm{exp})\pm 0.0007(\mathrm{theory})$$
(27)
The CLEO Collaboration Exp\_03 obtains from studies of the $`B^0\pi ^{}l^+\nu _l`$ branching fraction and $`q^2`$ distributions, using LCSR for $`0q^2<16\mathrm{GeV}^2`$ and lattice QCD for $`16\mathrm{GeV}^2q^2<q_{\mathrm{max}}^2`$,
$$|V_{ub}|_{\mathrm{CLEO}}=0.0032\pm 0.0003(\mathrm{exp})_{0.0004}^{+0.0006}(\mathrm{theory})$$
(28)
We see that both determinations of $`|V_{ub}|`$ are in an excellent agreement and that in both cases the error is dominated by uncertainties in the theoretical treatment. We have also calculated partially integrated branching ratios,
$$(q_1^2q^2<q_2^2)=\frac{_{\mathrm{exp}}^{\mathrm{total}}(B^0\pi ^{}l^+\nu _l)}{\mathrm{\Gamma }}_{q_1^2}^{q_2^2}𝑑q^2\frac{\mathrm{d}\mathrm{\Gamma }}{\mathrm{d}q^2}$$
(29)
Theoretical uncertainties partially cancel in the ratio $`_{q_1^2}^{q_2^2}𝑑q^2\frac{\mathrm{d}\mathrm{\Gamma }}{\mathrm{d}q^2}/\mathrm{\Gamma }`$. Our results are compiled in the Table 1. There it can be seen that they compare reasonably well with those quoted in Ref. Exp\_03 .
Finally, at each value of $`q^2`$ we also compute the spread of values obtained for the $`f^+`$ form factor when the five different quark-antiquark interactions are used. This procedure gives us a further theoretical error on $`f^+(q^2)`$ at fixed $`q^2`$ and by adding it in quadrature to that obtained previously from uncertainties on the $`(\left[g_{B^{}B\pi }f_B^{}\right]_{\mathrm{Latt}\mathrm{QCD}},f^+(0)_{\mathrm{LCSR}})`$ pair, we determine the theoretical error bands shown in Fig. 4.
## III Semileptonic $`D\pi `$ and $`DK`$ Decays
As a further test of our predictions for the $`B\pi `$ semileptonic process, we present results for the $`D\pi `$ and $`DK`$ decays for which there are precise experimental data pdg :
$`_{\mathrm{exp}}(D^0\pi ^{}e^+\nu _e)`$ $`=`$ $`(3.6\pm 0.6)\times 10^3,\mathrm{\Gamma }_{\mathrm{exp}}(D^0\pi ^{}e^+\nu _e)=\left(5.8\pm 1.0\right)\times 10^{12}\mathrm{MeV}`$ (30)
$`_{\mathrm{exp}}(D^0K^{}e^+\nu _e)`$ $`=`$ $`(3.58\pm 0.18)\times 10^2,\mathrm{\Gamma }_{\mathrm{exp}}(D^0K^{}e^+\nu _e)=\left(57\pm 3\right)\times 10^{12}\mathrm{MeV},`$ (31)
with life time $`\tau _{D^0}=(410.3\pm 1.5)\times 10^{15}\mathrm{s}`$ and $`|V_{cd}|=0.224\pm 0.003`$, $`|V_{cs}|=0.9737\pm 0.0007`$.
In the last two years there has been renewed interest in these decays. The first three-flavor lattice QCD results Latt\_D\_05 have appeared, superseding the old quenched ones Latt\_01 ; Latt\_01\_bis ; Latt\_D\_95 , and the BES Exp\_D\_04 and CLEO Exp\_D\_05 collaborations have new measurements of the branching ratios
$`\mathrm{BES}:(D^0\pi ^{}e^+\nu _e)`$ $`=`$ $`(3.3\pm 1.3)\times 10^3,(D^0K^{}e^+\nu _e)=(3.8\pm 0.5)\times 10^2`$ (32)
$`\mathrm{CLEO}:{\displaystyle \frac{(D^0\pi ^{}e^+\nu _e)}{(D^0K^{}e^+\nu _e)}}`$ $`=`$ $`0.082\pm 0.006\pm 0.005`$ (33)
Both collaborations have also determined the form factor at $`q^2=0`$
$`\mathrm{BES}:f_\pi ^+(0)`$ $`=`$ $`0.73\pm 0.14\pm 0.06,f_K^+(0)=0.78\pm 0.04\pm 0.03,{\displaystyle \frac{f_\pi ^+(0)}{f_K^+(0)}}=0.93\pm 0.19\pm 0.07`$ (34)
$`\mathrm{CLEO}:{\displaystyle \frac{f_\pi ^+(0)}{f_K^+(0)}}`$ $`=`$ $`0.86\pm 0.07_{0.04}^{+0.06}\pm 0.01`$ (35)
In the following we will apply the NRCQM developed for the $`B\pi `$ decay to the description of these $`D`$-meson semileptonic transitions. All formulae of Sect. II can be used here with the obvious replacements: $`BD,B^{}D^{}`$ for the $`D^0\pi ^{}l^+\nu _l`$ process, and $`BD`$, $`\pi K`$, $`B^{}D_s^{}`$ for the $`D^0K^{}l^+\nu _l`$ process. We will use $`m_{D^0}=1864.6\mathrm{MeV}`$, $`m_D^{}=2010\mathrm{MeV}`$, $`m_{D_s^{}}=2112.1\mathrm{MeV}`$ and $`m_K^{}=493.68\mathrm{MeV}`$.
### III.1 $`D\pi l\overline{\nu }_l`$
Since there is phase space for the $`D^{}D\pi `$ decay to occur, the $`g_{D^{}D\pi }`$ hadronic constant has been experimentally measured (CLEO gdstar )
$$g_{D^{}D\pi }=17.9\pm 0.3\pm 1.9$$
(36)
Taking $`f_D^{}=(234\pm 20)\mathrm{MeV}`$ from Ref. Latt\_fbstar , we find<sup>12</sup><sup>12</sup>12Note that the lattice QCD simulation of Ref. gdstar\_latt measured $`g_{D^{}D\pi }=18.8\pm 2.3\pm 2.0`$ in good agreement with Eq. (36).
$$\left[g_{D^{}D\pi }f_D^{}\right]_{\mathrm{Exp}\mathrm{Latt}}=4.2\pm 0.6\mathrm{GeV}$$
(37)
where we have added errors in quadrature. In Ref. NRCQM\_def and using the same set of NRCQM’s, we found a value of $`4.9\pm 0.5\mathrm{GeV}`$ for the above product, in reasonable agreement with Eq. (37). The value quoted in Eq. (37) determines the $`D^{}`$-pole contribution, above $`q^2=0`$, to $`f^+`$ and adding it to the valence quark contribution we obtain the results shown in Fig. 5. We find excellent agreement between our description of the form factor and that provided by the unquenched lattice simulation of Ref. Latt\_D\_05 . As can be seen in the figure, the $`D^{}`$-pole contribution is dominant above $`q^2=1.5\mathrm{GeV}^2`$ and it remains sizeable down to $`0.5\mathrm{GeV}^2`$. We do not see the need to Omnès improve the NRCQM description of the decay since it is quite good for the whole $`q^2`$ range. On the other hand, we see that the pion energy ranges from $`m_\pi `$ up to about $`1\mathrm{GeV}`$, which was also the maximum value for $`E_\pi `$ in the five NRCQM data-points used in the subtracted Omnès representation for the semileptonic $`B\pi `$ decay depicted in Fig. 4. This reinforces our belief in the reliability of our determination of $`|V_{ub}|`$ presented in Subsect. II.5.
Considering our theoretical uncertainties (errors on Eq. (37) and the spread of results obtained when different quark-antiquark interactions are considered) together with the experimental uncertainties on $`|V_{cd}|`$ quoted above, we find
$$\mathrm{\Gamma }_{\mathrm{this}\mathrm{work}}(D^0\pi ^{}e^+\nu _e)=[5.2\pm 0.1(\mathrm{exp}:|V_{cd}|)\pm 0.5(\mathrm{theory})]\times 10^{12}\mathrm{MeV}$$
(38)
in good agreement with Eq. (30). We also obtain
$$f_\pi ^+(0)=0.63\pm 0.02$$
(39)
compatible within errors with both the BES ($`0.73\pm 0.15`$) and the Fermilab-MILC-HPQCD ($`0.64\pm 0.07`$) results.
Finally in the left plot of Fig. 6, we compare the NRCQM predictions for the ratio $`f^+(q^2)/f^+(0)`$ with a pole form recently fitted to data by the FOCUS Collaboration focus .
### III.2 $`DKl\overline{\nu }_l`$
Since there is no phase space for the $`D_s^{}DK`$ decay, we will estimate the $`g_{D_s^{}DK}`$ coupling from the value quoted for $`g_{D^{}D\pi }`$ in Eq. (36). The parameter $`\widehat{g}`$ defined in Eq. (14) describes the strong coupling of charmed mesons as well as of beauty mesons to the members of the octet of light pseudoscalars. We will assume flavor SU(3) symmetry for this basic quantity in the heavy quark chiral effective theory, and thus we will use g\_02
$$g_{D_s^{}DK}\frac{2\widehat{g}\sqrt{m_Dm_{D_s^{}}}}{f_K}g_{D^{}D\pi }\frac{\sqrt{m_{D_s^{}}}}{\sqrt{m_D^{}}}\frac{f_\pi }{f_K}15.3\pm 1.6$$
(40)
where we have taken $`f_K/f_\pi 1.2`$ from Ref. Latt\_fK , and have kept some SU(3) flavor breaking terms in the masses of the charmed vector mesons and in the kaon decay constant. Taking $`f_{D_s^{}}=(254\pm 15)\mathrm{MeV}`$ from Ref. Latt\_fbstar , we find
$$\left[g_{D_s^{}DK}f_{D_s^{}}\right]_{\mathrm{SU}(3)\mathrm{Latt}}=3.9\pm 0.5\mathrm{GeV}$$
(41)
Our results for this decay are shown in Fig. 7. Several comments are in order:
1. We find reasonable agreement with the three-flavor lattice QCD results Latt\_D\_05 up to kaon energies of the order of $`1\mathrm{GeV}`$, which cover the whole $`q^2`$ range accessible in the physical decay. Discrepancies with lattice data are now more sizeable than in the $`D\pi `$ case and lattice unquenched data favor values for $`\left[g_{D_s^{}DK}f_{D_s^{}}\right]_{\mathrm{SU}(3)\mathrm{Latt}}`$ smaller than the one used in our calculation (see Eq. (40)). Theoretical errors for $`f^+`$ are, in this case, mostly due to the uncertainties on $`\left[g_{D_s^{}DK}f_{D_s^{}}\right]_{\mathrm{SU}(3)\mathrm{Latt}}`$. Nevertheless, we would like to point out that uncertainties on the value of $`g_{D_s^{}DK}`$ might be larger that those quoted in Eq. (40), since flavor SU(3) corrections to the relation $`g_{D_s^{}DK}2\widehat{g}\sqrt{m_Dm_{D_s^{}}}/f_K`$ could be large ($`m_s/m_cm_{d,u}/m_c`$).
2. The contribution of the vector resonance is less important than in the $`B\pi `$ and $`D\pi `$ decays, since the $`D_s^{}`$ is located relatively far from $`\sqrt{q_{\mathrm{max}}^2}`$.
3. Our predictions for $`f^+`$ at negative values of $`q^2`$, which do not enter into the phase space integral, suffer from larger uncertainties, since in that region the transferred momentum is larger than $`1\mathrm{GeV}`$ and, as for the $`B\pi `$ case, relativistic effects could became important. One could Omnès improve the NRCQM to achieve a better description of the form factor in the negative $`q^2`$ region.
In the right plot of Fig. 6, we compare the NRCQM predictions for the ratio $`f^+(q^2)/f^+(0)`$ with recently measured data from the FOCUS Collaboration focus and find satisfactory and reassuring agreement.
For the integrated width, we find
$$\mathrm{\Gamma }_{\mathrm{this}\mathrm{work}}(D^0K^{}e^+\nu _e)=\left[66\pm 3(\mathrm{theory})\right]\times 10^{12}\mathrm{MeV}$$
(42)
which is about two standard deviations higher than the value quoted in Eq. (31).
Eqs. (30) and (31) lead to (adding errors in quadratures)
$$\frac{(D^0\pi ^{}e^+\nu _e)}{(D^0K^{}e^+\nu _e)}=0.101\pm 0.017$$
(43)
which turns out to be a bit higher, though compatible within errors, than the recent CLEO determination quoted in Eq. (33). For this ratio of branching fractions we find
$$\frac{(D^0\pi ^{}e^+\nu _e)}{(D^0K^{}e^+\nu _e)}|_{\mathrm{this}\mathrm{work}}=0.079\pm 0.008$$
(44)
in excellent agreement with the CLEO measurement. We also find
$$f_K^+(0)=0.79\pm 0.01,\frac{f_\pi ^+(0)}{f_K^+(0)}=0.80\pm 0.03$$
(45)
which compare well to the recent experimental measurements in Eqs. (34) and (35).
## IV Concluding Remarks
We have shown the limitations of a valence quark model to describe the $`B\pi `$, $`D\pi `$ and $`DK`$ semileptonic decays. As a first correction, we have included in each case the heavy–light vector resonance pole contribution. For the semileptonic $`B\pi `$ decay, the inclusion of the $`B^{}`$ degree of freedom provides a realistic $`q^2`$ dependence of the relevant form factor, $`f^+`$, from $`q_{\mathrm{max}}^2`$ down to around $`18\mathrm{GeV}^2`$. We then use a multiply-subtracted Omnès dispersion relation, which considerably diminishes the form factor dependence on the elastic $`\pi B\pi B`$ scattering amplitudes at high energies, to combine LCSR results at $`q^2=0`$ with NRCQM predictions in the high $`q^2`$ region. As a result we have been able to predict the $`f^+`$ form factor for all $`q^2`$ values accessible in the physical decay. We have used a Monte Carlo procedure and analyzed the predictions of five different quark-antiquark interactions to determine theoretical error bands for form factors and the decay width. This has allowed us to extract from the measured branching fraction the value $`|V_{ub}|=0.0034\pm 0.0003(\mathrm{exp})\pm 0.0007(\mathrm{theory})`$ in excellent agreement with the CLEO Collaboration determination of Ref. Exp\_03 . For the $`D\pi `$ semileptonic decay we have found excellent agreement between our model calculation (valence quark plus $`D^{}`$-pole contributions) of $`f^+`$ and the one obtained by the unquenched lattice simulation of Ref. Latt\_D\_05 . We found no need to Omnès-improve our calculation in this case. Our results $`\mathrm{\Gamma }(D^0\pi ^{}e^+\nu _e)=[5.2\pm 0.1(\mathrm{exp}:|V_{cd}|)\pm 0.5(\mathrm{theory})]\times 10^{12}\mathrm{MeV}`$ and $`f_\pi ^+(0)=0.63\pm 0.02`$ are in good agreement with experimental data. Finally for the $`DK`$ semileptonic decay we find good agreement, in the physical region, between our model calculation (valence quark plus $`D_s^{}`$-pole contributions) of $`f^+`$ and lattice data, from UKQCD Latt\_D\_95 and Fermilab-MILC-HPQCD Latt\_D\_05 , and also with recent measurements from the FOCUS Collaboration focus . Again our results $`(D^0\pi ^{}e^+\nu _e)/(D^0K^{}e^+\nu _e)=0.079\pm 0.008`$, $`f_K^+(0)=0.79\pm 0.01`$ and $`f_\pi ^+(0)/f_K^+(0)=0.80\pm 0.03`$ are in good agreement with experimental determinations by the CLEO Exp\_D\_05 and BES Exp\_D\_04 Collaborations.
###### Acknowledgements.
This research was supported by DGI and FEDER funds, under contracts BFM2002-03218, BFM2003-00856 and FPA2004-05616, by the Junta de Andalucía and Junta de Castilla y León under contracts FQM0225 and SA104/04, and it is part of the EU integrated infrastructure initiative Hadron Physics Project under contract number RII3-CT-2004-506078. C. Albertus wishes to acknowledge a grant from Junta de Andalucía. J.M. Verde-Velasco acknowledges a grant (AP2003-4147) from the Spanish Ministerio de Educación y Ciencia. J.M. Flynn acknowledges PPARC grant PPA/G/O/2002/00468, the hospitality of the Universidad de Granada and the Institute for Nuclear Theory at the University of Washington, and thanks the Department of Energy for partial support during the completion of this work.
## Appendix A Multiply Subtracted Omnès Dispersion Relation
Let the form factor<sup>13</sup><sup>13</sup>13The discussion below can be trivially generalized to any scattering amplitude or form factor with definite total angular momentum and isospin quantum numbers. $`f^+(s)`$ be analytic on the complex $`s`$ plane (Mandelstam’s hypothesis Ma58 of maximum analyticity) except for a cut $`L[s_{\mathrm{th}}=(m_B+m_\pi )^2,+\mathrm{}[`$ along the real positive $`s`$ axis, as demanded by Watson’s theorem Wa55 . For real values $`s<s_{\mathrm{th}}`$ the form factor is real which implies that the values of the form factor above and below the cut are complex conjugates of each other: $`f^+(s+\mathrm{i}ϵ)=f^+(s\mathrm{i}ϵ)^{}`$. For $`ss_{\mathrm{th}}`$, the form factor has a discontinuity across the cut and develops an imaginary part $`f^+(s+\mathrm{i}ϵ)f^+(s\mathrm{i}ϵ)=2\mathrm{i}\mathrm{I}\mathrm{m}f(s+\mathrm{i}ϵ)`$. Cauchy’s theorem implies that $`f^+(s)`$ can be written as a dispersive integral along the cut and performing one subtraction at $`s_0<s_{\mathrm{th}}`$ one gets:
$$f^+(s)=f^+(s_0)+\frac{ss_0}{\pi }_{s_{\mathrm{th}}}^+\mathrm{}\frac{dx}{xs_0}\frac{\mathrm{Im}f^+(x)}{xs},sL,s_0<s_{\mathrm{th}}$$
(46)
Depending on the asymptotic behavior of $`f^+(s)`$ at the extremes of the cut $`L`$, more subtractions may be needed to make the integral convergent. For the time being, let us assume that one subtraction is sufficient. The well known Omnès solution for the above dispersive representation is Omnes :
$$O(s)=f^+(s_0)\mathrm{exp}\left\{\frac{ss_0}{\pi }_{s_{\mathrm{th}}}^+\mathrm{}\frac{dx}{xs_0}\frac{\delta (x)}{xs}\right\},sL,s_0<s_{\mathrm{th}}$$
(47)
with $`\delta (s)`$ the elastic $`\pi B\pi B`$ phase shift<sup>14</sup><sup>14</sup>14Obviously $`\delta (s)`$ has to be defined as a continuous function of $`s`$. in the $`J^P=1^{}`$ and isospin $`1/2`$ channel (see Eq. (20)). $`O(s)`$ gives the physical form-factor since,
1. For $`ss_{\mathrm{th}}`$, we have
$`O(s\pm \mathrm{i}ϵ)`$ $`=`$ $`f^+(s_0)\mathrm{exp}\left\{{\displaystyle \frac{ss_0}{\pi }}\left[𝒫{\displaystyle _{s_{\mathrm{th}}}^+\mathrm{}}{\displaystyle \frac{dx}{xs_0}}{\displaystyle \frac{\delta (x)}{xs}}\pm \mathrm{i}\pi {\displaystyle \frac{\delta (s)}{ss_0}}\right]\right\}`$ (48)
$`=`$ $`e^{\pm \mathrm{i}\delta (s)}\left[f^+(s_0)\mathrm{exp}\left\{{\displaystyle \frac{ss_0}{\pi }}𝒫{\displaystyle _{s_{\mathrm{th}}}^+\mathrm{}}{\displaystyle \frac{dx}{xs_0}}{\displaystyle \frac{\delta (x)}{xs}}\right\}\right]`$
where $`𝒫`$ stands for principal part of the integral. Thus we have $`O(s+\mathrm{i}ϵ)=O(s\mathrm{i}ϵ)^{}`$, the function $`O`$ is real for $`s<s_{\mathrm{th}}`$ and has neither poles nor cuts, except for that required by Watson’s theorem: $`L[s_{\mathrm{th}},+\mathrm{}[`$. The discontinuity across this cut is given by $`O(s+\mathrm{i}ϵ)O(s\mathrm{i}ϵ)=2\mathrm{i}\mathrm{I}\mathrm{m}O(s+\mathrm{i}ϵ)`$ and by construction $`O(s_0)=f^+(s_0)`$. Thus, both $`f^+(s)`$ and $`O(s)`$ satisfy the same dispersion relation (Eq. (46)) and therefore both functions can differ at most by a polynomial with real coefficients, which should vanish at $`s=s_0`$. But, this polynomial is zero since:
2. The function $`O(s)`$ satisfies Watson’s theorem:
$$\frac{O(s+\mathrm{i}ϵ)}{O(s\mathrm{i}ϵ)}=e^{2\mathrm{i}\delta (s)}=\frac{f^+(s+\mathrm{i}ϵ)}{f^+(s\mathrm{i}ϵ)},s>s_{\mathrm{th}}$$
(49)
Performing $`n+1`$ subtractions, one can can produce a rank $`n+2`$ polynomial in the denominator of the dispersive integral of Eq. (46). Indeed, for $`sL`$
$$f^+(s)=P_n(s)+\frac{(ss_0)(ss_1)\mathrm{}(ss_n)}{\pi }_{s_{\mathrm{th}}}^+\mathrm{}\frac{dx}{(xs_0)(xs_1)\mathrm{}(xs_n)}\frac{\mathrm{Im}f^+(x)}{xs},s_0,\mathrm{},s_n<s_{\mathrm{th}}$$
(50)
with the rank $`n`$ polynomial $`P_n(s)`$ determined by the $`n+1`$ equations $`P_n(s_i)=f^+(s_i),i=0,1,\mathrm{},n`$,
$$P_n(s)=\underset{j=0}{\overset{n}{}}\alpha _j(s)f^+(s_j),\alpha _j(s)=\left[\underset{jk=0}{\overset{n}{}}\frac{ss_k}{s_js_k}\right]$$
(51)
Note that $`\alpha _j(s)`$ are rank $`n`$ polynomials, which satisfy
$$\underset{j=0}{\overset{n}{}}\alpha _j(s)=1$$
(52)
On the other hand for $`s>s_{\mathrm{th}}`$, we have from Eq. (20)
$$\mathrm{log}f^+(s+\mathrm{i}ϵ)\mathrm{log}f^+(s\mathrm{i}ϵ)=\mathrm{log}\frac{f^+(s+\mathrm{i}ϵ)}{f^+(s\mathrm{i}ϵ)}=2\mathrm{i}\delta (s)=2\mathrm{i}\mathrm{I}\mathrm{m}\left[\mathrm{log}f^+(s+\mathrm{i}ϵ)\right]$$
(53)
Thus in analogy to Eq. (50), assuming that the form factor does not vanish in $`𝒞\{s_{\mathrm{th}}\}`$<sup>15</sup><sup>15</sup>15Note that, we have already treated $`s=s_{\mathrm{th}}`$ as a branch point., or neglecting the contribution from the log cut if it has a finite branch point different from $`s_{\mathrm{th}}`$, we can write
$$\mathrm{log}f^+(s)=\widehat{P}_n(s)+\frac{(ss_0)(ss_1)\mathrm{}(ss_n)}{\pi }_{s_{\mathrm{th}}}^+\mathrm{}\frac{dx}{(xs_0)(xs_1)\mathrm{}(xs_n)}\frac{\delta (x)}{xs},sL$$
(54)
with
$$\widehat{P}_n(s)=\underset{j=0}{\overset{n}{}}\alpha _j(s)\mathrm{log}f^+(s_j)$$
(55)
From the above equation one readily finds the $`(n+1)`$-subtracted Omnès representation given in Eq. (19).
Finally, we would like to draw the reader’s attention to a subtle point. In Eq. (46) we have assumed that $`f^+`$ has no poles. However we know that if the scattering amplitude has a pole at $`s_R=M_R^2\mathrm{i}M_R\mathrm{\Gamma }_R`$ on its second Riemann sheet (resonance) or on the physical sheet (bound state with $`\mathrm{\Gamma }_R=0^+`$ and $`M_R^2<s_{\mathrm{th}}`$.), it might show up as a pole in the complex plane of $`f^+`$ (see Eq. (15)). On the other hand, the $`S`$-matrix depends on $`\mathrm{exp}(2\mathrm{i}\delta )`$, and thus one has the freedom to add factors of $`m\pi `$, for $`m`$ an integer, to the phase-shift without modifying the $`S`$-matrix. However, the Omnès representation of the form factor will definitely depend on the specific value chosen for the integer $`m`$.
To fix this ambiguity, we will assume that at threshold the phase shift should be $`\delta (s_{\mathrm{th}})=n_b\pi `$, where $`n_b`$ is the number of bound states in the channel, while $`\delta (\mathrm{})=k\pi `$, where $`k`$ is the number of zeros of the scattering amplitude on the physical sheet (this is Levinson’s theorem MS70 ). This choice for the phase shifts also takes into account the existence of poles in the scattering matrix. We demonstrate with a simple example in which a $`p`$-wave $`T`$-matrix is proportional to $`(ss_{\mathrm{th}})/(sM_R^2+\mathrm{i}M_R\mathrm{\Gamma }_R)`$. The phase shift is given by
$$\delta (s)=\pi +\mathrm{Arctan}\left[\frac{M_R\mathrm{\Gamma }_R}{sM_R^2}\right],s>s_{\mathrm{th}}$$
(56)
with $`\mathrm{Arctan}[\pi ,\pi [`$. This satisfies $`\delta (\mathrm{})=\pi `$ and, if $`M_R\mathrm{\Gamma }_R|sM_R^2|`$, it also leads to $`\delta (s_{\mathrm{th}})=\pi `$ or $`0`$ for a bound state or resonance respectively, in accordance with Levinson’s theorem. For simplicity, let us also assume $`\mathrm{\Gamma }_RM_R`$. In this circumstance we can approximate
$$\delta (s)\pi \left[1H(M_R^2s)\right]=\pi H(sM_R^2)$$
(57)
where $`H()`$ is the step function. Since
$$\frac{ss_0}{\pi }_{s_{\mathrm{th}}}^+\mathrm{}\frac{dx}{xs_0}\frac{\delta (x)}{xs}(ss_0)_{\mathrm{Max}(s_{\mathrm{th}},M_R^2)}^+\mathrm{}\frac{dx}{xs_0}\frac{1}{xs}=\mathrm{log}\left\{\frac{\mathrm{Max}(s_{\mathrm{th}},M_R^2)s_0}{\mathrm{Max}(s_{\mathrm{th}},M_R^2)s}\right\},$$
(58)
we find that the Omnès solution from a once-subtracted dispersion relation, Eq. (47), reads
$$O(s)f^+(s_0)\frac{\mathrm{Max}(s_{\mathrm{th}},M_R^2)s_0}{\mathrm{Max}(s_{\mathrm{th}},M_R^2)s}$$
(59)
It has a pole at $`s=M_R^2`$ for a resonance or at $`s=s_{\mathrm{th}}`$ in the case of a bound state. For the case of a resonance with a finite width, the pole will move to $`s=M_R^2+\mathrm{i}M_R\mathrm{\Gamma }_R`$. On the other hand, for a bound state, going beyond the approximation $`\delta (s)\pi `$ (see Eq. (21)), the form factor will be sensitive to the exact position of the pole ($`s=M_R^2`$), since the effective range parameters (scattering volume, …) will depend on $`M_R`$.
These conclusions can easily be generalized when a multiply-subtracted Omnès dispersion relation is used. |
warning/0506/cond-mat0506077.html | ar5iv | text | # Quantum mechanics as an approximation of statistical mechanics for classical fields
## 1 Introduction
The problem of coupling the quantum probabilistic model (which is based on the Hilbert space calculus) and the classical probabilistic model (which is based on the measure-theoretic calculus) was intensively discussed already by fathers of quantum mechanics, see, e.g., the correspondence between Einstein and Schrödinger . This is the problem of huge complexity, see, e.g., for different viewpoints and debates.
Now days there is a rather common opinion that the probabilistic structure of quantum mechanics cannot be considered simply as a special mathematical representation of classical (measure-theoretic) probability theory. Such an opinion is based merely on a number of no-go theorems, see appendix 2 for details. On the other hand, there are known various prequantum models that reproduce (at least some) features of quantum mechanics: De Broglie’s double solution theory, Bohmian mechanics, Nelson’s stochastic mechanics, t’ Hooft’s deterministic discrete models, see . All such models are either nonlocal (as Bohmian mechanics and Nelson’s stochastic mechanics) or reproduce only some (not all) predictions of quantum mechanics. In any event they do not contradict no-go theorems.
In author’s papers , there was proposed a new prequantum model: Prequantum Classical Statistical Field Theory, PCSFT. This model is very natural, because this is nothing else than standard classical statistical mechanics: a) state space is phase space; b) variables are functions on phase space; c) statistical states (describing ensembles of systems) are probability measures. One important point is that phase space is infinite-dimensional. Prequantum states (“hidden variables”) can be represented as vector fields, $`\psi (x)=(q(x),p(x)).`$ However, this is merely a technical mathematical feature of PCSFT. The crucial point is a tricky way in which PCSFT is projected onto quantum mechanics. This projection is asymptotic. Quantum mechanics can be considered as an approximative theory: the quantum average given by the von Neumann trace formula appears as the first nontrivial contribution into the classical average (given by the Lebesgue integral).
To clarify the main distinguishing features of our theory, PCSFT, we shall divide its presentation into a few steps. First we consider the finite dimensional case. Here all computations are reduced to asymptotic expansions of simple Gaussian integrals. Then we shall proceed to the infinite-dimensional case where we should consider Gaussian integrals over functional spaces. We also start with consideration of a toy-model of quantum mechanics over reals, see also , . This will simplify the understanding of PCSFT, because the prequantum phase-space model inducing the quantum model over complex numbers has a rather nontrivial geometric structure. And its presentation should be separated from asymptotic expansion of Gaussian integrals.
We remark that construction of a prequantum classical statistical model for finite-dimensional quantum mechanics is interesting not only from purely mathematical viewpoint. Since electron-spin and photon-polarization can be described by two dimensional complex spaces, in the finite-dimensional case our approach shows that that its is possible to construct a pure classical phase-space models for these quantum phenomena, see appendix 1.
In our approach quantum mechanics is an approximative theory. It predicts statistical averages only with some precision. In principle, there might be found deviations of averages calculated within the quantum formalism from experimental averages (which are supposed to be equal to classical averages given by our model). But at the moment our predictions is not of a high value for experimentalists, because PCSFT does not predict the magnitude of a small parameter $`\alpha `$ in the asymptotic representation of classical averages. In , we speculated that the small parameter of PCSFT $`\alpha `$ (the dispersion of prequantum fluctuations) can be chosen equal to the Planck constant $`\mathrm{}.`$ But that speculation was not justified. We notice that our asymptotic considerations are totally different from the standard considerations on the classical limit of quantum mechanics: obtaining classical phase space mechanics as the limit of quantum mechanics when the Planck constant $`\mathrm{}`$ (which is considered as a small parameter) goes to zero. In our approach when the small parameter $`\alpha `$ goes to zero we obtain quantum theory as the limit case of classical (and not vice versa). In particular, neglecting by $`\mathrm{}`$ induces neglecting by “quantum rotation”, spin. But in PCSFT such degrees of freedom are not neglected, see appendix 1 for details.
## 2 The Taylor approximation of averages for functions of random variables
Here we follow chapter 11 of the book of Elena Ventzel. This book was written in the form of precise instructions which student should follow to solve a problem:
“In practice we have very often situations in that, although investigated function of random arguments is not strictly linear, but it differs practically so negligibly from a linear function that it can be approximately considered as linear. This is a consequence of the fact that in many problems fluctuations of random variables play the role of small deviations from the basic law. Since such deviations are relatively small, functions which are not linear in the whole range of variation of their arguments are almost linear in a restricted range of their random changes,” , p. 238.
Let $`y=f(x).`$ Here in general $`f`$ is not linear, but it does not differ so much from linear on some interval $`[m_x\delta ,m_x+\delta ],`$ where $`x=x(\omega )`$ is a random variable and
$$m_xEx=x(\omega )𝑑𝐏(\omega )$$
is its average. Here $`\delta >0`$ is sufficiently small. Student of a military college should approximate $`f`$ by using the first order Taylor expansion at the point $`m_x:`$
$$y(\omega )f(m_x)+f^{}(m_x)(x(\omega )m_x).$$
(1)
By taking the average of both sides he obtains:
$$m_yf(m_x).$$
(2)
The crucial point is that the linear term $`f^{}(m_x)(x(\omega )m_x)`$ does not give any contribution! Further Elena Ventzel pointed out , p. 245: “For some problems the above linearization procedure may be unjustified, because the method of linearization may be not produce a sufficiently good approximation. In such cases to test the applicability of the linearization method and to improve results there can be applied the method which is based on preserving not only the linear term in the expansion of function, but also some terms of higher orders.”
Let $`y=f(x).`$ Student now should preserve the first three terms in the expansion of $`f`$ into the Taylor series at the point $`m_x:`$
$$y(\omega )f(m_x)+f^{}(m_x)(x(\omega )m_x)+\frac{1}{2}f^{\prime \prime }(m_x)(x(\omega )m_x)^2.$$
(3)
Hence
$$m_yf(m_x)+\frac{\sigma _x^2}{2}f^{\prime \prime }(m_x),$$
(4)
where
$$\sigma _x^2=E(xm_x)^2=(x(\omega )m_x)^2𝑑𝐏(\omega )$$
is the dispersion of the random variable $`x.`$
Let us now consider the special case of symmetric fluctuations:
$$m_x=0$$
and let us restrict considerations to functions $`f`$ such that
$$f(0)=0.$$
Then we obtain the following special form of (4):
$$m_y\frac{\sigma _x^2}{2}f^{\prime \prime }(0).$$
(5)
We emphasize again that the first derivative does not give any contribution into the average.
Thus at the some level of approximation we can calculate averages not by using the Lebesgue integral (as we do in classical probability theory), but by finding the second derivative. Such a “calculus of probability” would match well with experiment. I hope that reader has already found analogy with the quantum calculus of probabilities. But for a better expression of this analogy we shall consider the multi-dimensional case. Let now $`x=(x_1,\mathrm{},x_n),`$ so we consider a system of $`n`$ random variables. We consider the vector average: $`m_x=(m_{x_1},\mathrm{},m_{x_n})`$ and the covariance matrix:
$$B_x=(B_x^{ij}),B_x^{ij}=E(x_im_{x_i})(x_jm_{x_j}).$$
We now consider the random variable
$$y(\omega )=f(x_1(\omega ),\mathrm{},x_n(\omega )).$$
By using the Taylor expansion we would like to obtain an algorithm for approximation of the average $`m_y.`$ We start directly from the second order Taylor expansion:
$$y(\omega )f(m_{x_1},\mathrm{},m_{x_n})+\underset{i=1}{\overset{n}{}}\frac{f}{x_i}(m_{x_1},\mathrm{},m_{x_n})(x_i(\omega )m_{x_i})$$
$$+\frac{1}{2}\underset{i,j=1}{\overset{n}{}}\frac{^2f}{x_ix_j}(m_{x_1},\mathrm{},m_{x_n})(x_i(\omega )m_{x_i})(x_j(\omega )m_{x_j}),$$
(6)
and hence:
$$m_yf(m_{x_1},\mathrm{},m_{x_1})+\frac{1}{2}\underset{i,j=1}{\overset{n}{}}\frac{^2f}{x_ix_j}(m_{x_1},\mathrm{},m_{x_1})B_x^{ij}.$$
(7)
By using the vector notations we can rewrite the previous formulas as:
$$y(\omega )f(m_x)+(f^{}(m_x),x(\omega )m_x)+\frac{1}{2}(f^{\prime \prime }(m_x)(x(\omega )m_x),x(\omega )m_x).$$
(8)
and
$$m_yf(m_x)+\frac{1}{2}\mathrm{Tr}\mathrm{B}_\mathrm{x}\mathrm{f}^{\prime \prime }(\mathrm{m}_\mathrm{x}).$$
(9)
Let us again consider the special case: $`m_x=0`$ and $`f(0)=0.`$ We have:
$$m_y\frac{1}{2}\mathrm{Tr}\mathrm{B}_\mathrm{x}\mathrm{f}^{\prime \prime }(0).$$
(10)
We now remark that the Hessian $`f^{\prime \prime }(0)`$ is always a symmetric operator. Let us now represent $`f`$ by its second derivative at zero:
$$fA=f^{\prime \prime }(0).$$
Then we see that, at some level of approximation, instead of operation with Lebesgue integrals, one can use linear algebra:
$$m_y\frac{1}{2}\mathrm{Tr}\mathrm{B}_\mathrm{x}\mathrm{A}$$
(11)
We now proceed in mathematically rigorous way, namely, we shall estimate the reminder which was neglected in the approximative formula for average. We also formalize correspondence between classical and quantum statistical models.
## 3 Classical and quantum statistical models
### 3.1 Classical statistical model
Classical statistical mechanics on phase space $`\mathrm{\Omega }_{2n}=𝐑^n\times 𝐑^n`$ can be considered as a special classical statistical model. In general a classical statistical model is defined in the following way:
a). Physical states $`\omega `$ are represented by points of some set $`\mathrm{\Omega }`$ (state space).
b). Physical variables are represented by functions $`f:\mathrm{\Omega }𝐑`$ belonging to some functional space $`V(\mathrm{\Omega }).`$<sup>1</sup><sup>1</sup>1The choice of a concrete functional space $`V(\mathrm{\Omega })`$ depends on various physical and mathematical factors.
c). Statistical states are represented by probability measures on $`\mathrm{\Omega }`$ belonging to some class $`S(\mathrm{\Omega }).`$
d). The average of a physical variable (which is represented by a function $`fV(\mathrm{\Omega }))`$ with respect to a statistical state (which is represented by a probability measure $`\rho S(\mathrm{\Omega }))`$ is given by
$$<f>_\rho _\mathrm{\Omega }f(\omega )𝑑\rho (\omega ).$$
(12)
A classical statistical model is a pair
$$M=(S(\mathrm{\Omega }),V(\mathrm{\Omega })).$$
In classical statistical mechanics $`\mathrm{\Omega }=\mathrm{\Omega }_{2n}`$ is phase space, $`V(\mathrm{\Omega }_{2n})=C^{\mathrm{}}(\mathrm{\Omega }_{2n})`$ is the space of all smooth functions on phase space, $`S(\mathrm{\Omega }_{2n})`$ is the space $`PM(\mathrm{\Omega }_{2n})`$ of all probability measures on phase space and the average is given by the Lebesgue integral on the $`\sigma `$-algebra of Borel subsets of $`\mathrm{\Omega }_{2n}.`$
Remark 3.1. We emphasize that the space of variables $`V(\mathrm{\Omega })`$ need not coincide with the space of all random variables $`RV(\mathrm{\Omega })`$ – measurable functions $`\xi :\mathrm{\Omega }𝐑.`$ For example, if $`\mathrm{\Omega }`$ is a differentiable manifold, it is natural to choose $`V(\mathrm{\Omega })`$ consisting of smooth functions; if $`\mathrm{\Omega }`$ is an analytic manifold, it is natural to choose $`V(\mathrm{\Omega })`$ consisting of analytic functions and so on. The space of statistical states $`S(\mathrm{\Omega })`$ need not coincide with the space of all probability measures $`PM(\mathrm{\Omega }).`$ For example, for some statistical model $`S(\mathrm{\Omega })`$ may consist of Gaussian measures.
### 3.2 Real finite-dimensional quantum mechanics
We shall use a toy model of quantum mechanics which based on the real space. Statistical features of the correspondence between a prequantum classical statistical model and quantum mechanics are more evident for this toy model. Denote the algebra of all $`(m\times m)`$ real matrices by the symbol $`M^{(r)}(m).`$ We denote by $`𝒟^{(r)}(m)`$ the class of nonnegative symmetric trace-one matrices $`\rho M^{(r)}(m).`$ We call them “density operators.” We denote by $`_s^{(r)}(m)`$ the class of all symmetric matrices. In the quantum model (for the $`m`$-dimensional real space) statistical states (describing ensembles of systems prepared for measurement) are represented by density matrices and quantum observables by matrices belonging $`_s^{(r)}(m).`$ The quantum average of an observable $`A_s^{(r)}(m)`$ with respect to a statistical state $`\rho 𝒟^{(r)}(m)`$ is given by the von Neumann trace class formula :
$$<A>_\rho =\mathrm{Tr}\rho \mathrm{A}.$$
(13)
In the operator representation observables and density matrices are corresponding classes of $`𝐑`$-linear operators. Denote the quantum model by
$$N_{\mathrm{quant}}^{(r)}=(𝒟^{(r)}(m),_s^{(r)}(m)).$$
If $`m=1,`$ then quantum observables are given by real numbers (operators of multiplication by real numbers on the real line) and there is only one statistical state $`\rho =1.`$ Here $`<A>_\rho =\rho A=A.`$
### 3.3 Complex finite-dimensional quantum mechanics
Denote the algebra of all $`(m\times m)`$ complex matrices by the symbol $`M^{(c)}(m).`$ We denote by $`𝒟^{(c)}(m)`$ the class of nonnegative symmetric trace-one matrices $`\rho M^{(c)}(m).`$ We call them “density operators.” We denote by $`_s^{(c)}(m)`$ the class of all symmetric matrices. In the quantum model (for the $`m`$-dimensional complex space) statistical states (describing ensembles of systems prepared for measurement) are represented by density matrices and quantum observables by matrices belonging $`_s^{(c)}(m).`$ The quantum average is given by (13). In the operator representation observables and density matrices are corresponding classes of $`𝐂`$-linear operators. Denote the quantum model by
$$N_{\mathrm{quant}}^{(c)}=(𝒟^{(c)}(m),_s^{(c)}(m)).$$
If $`m=1,`$ then quantum observables are given by real numbers (operators of multiplication by real numbers on the complex plane) and and there is only one statistical state $`\rho =1.`$ Here $`<A>_\rho =\rho A=A.`$
### 3.4 Complex quantum mechanics
Denote by $`H_c`$ a complex (separable and infinite-dimensional) Hilbert space. Denote the algebra of all bounded operators $`A:H_cH_c`$ by the symbol $`(H_c).`$ The real linear subspace of $`(H_c)`$ consisting of self-adjoint operators is denoted by the symbol $`_s(H_c).`$ We denote by $`𝒟(H_c)`$ the class of nonnegative trace-one operators $`\rho _s(H_c).`$ These are von Neumann density operators. In the quantum model statistical states (describing ensembles<sup>2</sup><sup>2</sup>2We follow so called ensemble interpretation of quantum mechanics: Einstein, Margenau, Ballentine, Balian, Nieuwenhuizen and many others, see, e.g., . By such an interpretation even a pure state (normalized vector of $`H_c)`$ represents an ensemble. In the orthodox Copenhagen interpretation a pure state represents the state of an individual system, e.g., electron. of systems prepared for measurement) are represented by density operators and quantum observables by operators from $`_s(H_c).`$ The quantum average is given by (13). Denote the quantum model by
$$N_{\mathrm{quant}}(H_c)=(𝒟(H_c),_s(H_c)).$$
## 4 Taylor approximation of classical averages: one dimensional case
States of systems are represented by real numbers, $`qQ=𝐑.`$ Ensembles of such systems are described by probability measures on the real line, statistical states. We consider a special class of preparation procedures. They produce ensembles of systems described by Gaussian probability distributions on $`Q`$ having the zero mean value and dispersion
$$\sigma ^2(\mu )=\alpha +O(\alpha ^2),$$
(14)
where as always $`|O(\alpha ^2)|C\alpha ^2`$ for some constant $`C`$ and a sufficiently small $`\alpha .`$ The crucial point is that $`\alpha `$ is a small parameter of our model. Denote this class of probability distributions by the symbol $`S_G^\alpha (Q).`$
For a probability $`\mu S_G^\alpha (Q),`$ we have:
$$d\mu (q)=\frac{e^{\frac{q^2}{2(\alpha +O(\alpha ^2))}}dq}{\sqrt{2\pi (\alpha +O(\alpha ^2))}}.$$
(15)
We recall that, for a probability with the zero mean value, its dispersion is given by
$$\sigma ^2(\mu )=\frac{1}{\sqrt{2\pi (\alpha +O(\alpha ^2))}}_{\mathrm{}}^{\mathrm{}}q^2e^{\frac{q^2}{2(\alpha +O(\alpha ^2))}}𝑑q.$$
(16)
As was already pointed out, we consider $`\alpha `$ as a small parameter. Therefore Gaussian probability distributions are very sharply concentrated around the point $`q_0=0.`$ By using the terminology of functional analysis we say that $`\{\mu \mu (\alpha )\}`$ is a $`\delta `$-family: $`lim_{\alpha 0}\mu (\alpha )=\delta `$ in the sense of theory of distributions.
In the approximation $`\alpha =0`$ all systems are located at a single point, namely, $`q_0.`$ However, a finer description (in that $`\alpha `$ can not be neglected) provides the picture of Gaussian bells concentrated nearby $`q_0.`$ We remark that in average a system cannot go far away from $`q_0.`$ By using the Chebyshov inequality one obtain for any $`C>0:`$
$$\mu \{q:|q|>C\}\frac{\alpha +O(\alpha ^2)}{C^2}O,\alpha 0.$$
(17)
But the probabilistic inequality (17) does not exclude the possibility that some system could move far from $`q_0`$ (of course, with a small probability).
We also introduce a class of physical variables in the classical statistical model under consideration:
a) $`fC^{\mathrm{}}(𝐑),`$ a smooth function;
b) $`f(0)=0;`$
c) $`|f^{(4)}(q)|c_fe^{r_f|q|},c_f,r_f0.`$
Denote this functional space by the symbol $`𝒱(Q),Q=𝐑.`$
The restriction to the growth of the fourth derivative will be used when we shall consider the Taylor expansion of $`f`$ up two the fourth term. The exponential growth implies integrability with respect to any Gaussian measure.
Lemma 4.1. Let $`fC^n`$ (so it is $`n`$ times continuously differentiable) and let its $`n`$th derivative has the exponential growth. Then all derivatives of orders $`n=0,\mathrm{},n1,`$ also have the exponential growth (in particular, $`f(q)`$ grows exponentially).
Proof. Under these conditions we can use the Taylor expansion with the integral remainder:
$$f(q)=f(0)+f^{}(0)q+\frac{f^{\prime \prime }(0)q^2}{2}+\frac{f^{(3)}(0)q^3}{3!}+\mathrm{}+\frac{q^n}{n!}_0^1(1\theta )^{n1}f^{(n)}(\theta q)𝑑\theta .$$
(18)
Since the growth of any polynomial can be compensated by decreasing of the $`e^{r|q|},`$ by using the exponential estimate for the $`n`$th derivative we obtain:
$$|f(q)|=C_1e^{r|q|}+C_2\frac{q^n}{n!}_0^1(1\theta )^{n1}e^{r|q\theta |}𝑑\theta Ce^{r|q|}.$$
(19)
Here all constants depend on $`f.`$
This simple exercise from the course of analysis will be useful in our further considerations. We defined the following classical statistical model on the real line:
A). States of systems are real numbers.
B). Statistical states (ensembles of systems) are represented by Gaussian probabilities having zero average and dispersion $`\sigma ^2(\mu )=\alpha +O(\alpha ^2),\alpha 0.`$
C). Physical variables are smooth functions with exponentially growing fourth derivative which map zero into itself.
We denote this model by $`N_{\mathrm{class}}^\alpha =(S_G^\alpha (Q),𝒱(Q)).`$ As always in classical statistical physics, the average of a physical variable $`f𝒱(Q)`$ with respect to an ensemble of systems which is described by a probability $`\mu S_G^\alpha (Q)`$ is given by the integral:
$$<f>_\mu =\frac{1}{\sqrt{2\pi (\alpha +O(\alpha ^2))}}_{\mathrm{}}^{\mathrm{}}f(q)e^{\frac{q^2}{2(\alpha +O(\alpha ^2))}}𝑑q.$$
(20)
Since $`\alpha `$ is a parameter of the model, we can consider averages as functions of $`\alpha :<f>_\mu <f>_\mu (\alpha ).`$ We are interested in the asymptotic expansion of averages when $`\alpha 0.`$ In particular, such an asymptotic expansion will give us the possibility to calculate averages approximately.
Lemma 4.2. Let $`f𝒱(Q)`$ and let $`\mu S_G^\alpha (Q).`$ Then
$$<f>_\mu (\alpha )=\frac{\alpha }{2}f^{\prime \prime }(0)+O(\alpha ^2).$$
(21)
Proof. We start with the scaling of the state variable:
$$q=\sigma (\mu )x$$
(22)
We have:
$$<f>_\mu (\alpha )=\frac{1}{\sqrt{2\pi }}_{\mathrm{}}^{\mathrm{}}f(\sigma (\mu )x)e^{\frac{x^2}{2}}𝑑x.$$
(23)
We now expand $`f(\sigma (\mu )x)`$ by using the fourth order Taylor formula with the integral remainder, see Lemma 4.1:
$$<f>_\mu (\alpha )=\frac{\sigma ^2(\mu )}{2}f^{\prime \prime }(0)$$
(24)
$$+\frac{\sigma ^4(\mu )}{4!\sqrt{2\pi }}_{\mathrm{}}^{\mathrm{}}x^4\left(_0^1(1\theta )^3f^{(4)}(\sigma (\mu )x\theta )𝑑\theta \right)e^{\frac{x^2}{2}}𝑑x.$$
We recall that for a Gaussian measure with zero mean value all odd momenta are equal to zero. This is an important point of our considerations. This imply that the first nonzero contribution to the classical average is given by the second derivative – quadratic term. Disappearance of the third order term implies the asymptotics $`O(\alpha ^2).`$ We now estimate the remainder to obtain this asymptotics:
$$|R(f,\mu )|\frac{C\sigma ^4(\mu )}{4!\sqrt{2\pi }}_{\mathrm{}}^{\mathrm{}}x^4\left(_0^1(1\theta )^3e^{r\sigma (\mu )|x|\theta }𝑑\theta \right)e^{\frac{x^2}{2}}𝑑x.$$
Since we consider $`\alpha `$ as a small parameter, we can assume that $`|\sigma (\mu )|1`$ in the exponential function. Thus:
$$|R(f,\mu )|C^{}\sigma ^4(\mu )_{\mathrm{}}^{\mathrm{}}x^4e^{r|x|\frac{x^2}{2}}𝑑x.$$
Since $`\sigma ^2(\mu )=\alpha +O(\alpha ^2),`$ we have that $`R(f,\mu )=O(\alpha ^2),\alpha 0.`$
We consider the dispersion $`\sigma ^2(\mu )`$ as the intensity of fluctuations in the ensemble of systems. We define the relative average with respect to the intensity of fluctuations by normalizing the average by the main term – namely, $`\alpha `$ – in the intensity of fluctuations:
$$f_\mu =\frac{<f>_\mu }{\alpha }.$$
Of course, $`f_\mu `$ is also a function of the parameter $`\alpha :`$
$$f_\mu (\alpha )=\frac{<f>_\mu (\alpha )}{\alpha }.$$
Corollary 4.1. Let $`f𝒱(Q)`$ and let $`\mu S_G^\alpha (Q).`$ Then
$$f_\mu =\frac{f^{\prime \prime }(0)}{2}+O(\alpha ).$$
(25)
In particular,
$$\underset{\alpha 0}{lim}f_\mu (\alpha )=\frac{f^{\prime \prime }(0)}{2}.$$
(26)
Proposition 4.1. We have:
$$f_\mu =\frac{<f>_\mu }{\sigma ^2(\mu )}+O(\alpha ).$$
(27)
Remark 4.1 (About $`1/2)`$ The second term in the Taylor formula gives the factor 1/2 which looks rather bothering in our asymptotic formula for the classical average. This factor will disappear in the complex representation and the formula will become nicer.
We have shown that $`\frac{f^{\prime \prime }(0)}{2}`$ gives the approximation of the (classical) relative average. The precision of such an approximation is $`\alpha .`$ If the level of development of measurement technology is such that all contributions of the magnitude $`\alpha `$ are neglected in measurements, then averages can be calculated by using the following simple rule:
$$f_\mu ^{\mathrm{approx}}=\left[\frac{<f>_\mu }{\sigma ^2(\mu )}\right]^{\mathrm{approx}}=\frac{f^{\prime \prime }(0)}{2}.$$
(28)
At the first sight such averages have nothing to do with classical averages given by integrals. There could be even presented an interpretation of physics claiming that rules of classical probability theory are violated and relating the exotic rule (28) for calculating of averages to special features of systems under consideration (and not to a special approximation procedure for averages).
Finally, we remark that calculation of averages by (28) is essentially simpler than classical probabilistic averages given by Lebesgue integrals.
## 5 Taylor approximation of classical averages: multidimensional case
States are vectors $`qQ=𝐑^m;`$ statistical states are Gaussian distributions with the zero mean value and the dispersion $`\sigma ^2(\mu )=\alpha +O(\alpha ^2).`$ Denote this class of probabilities by the symbol $`S_G^\alpha (Q).`$ We introduce the scalar product and norm on $`Q:`$
$$(\xi ,q)=\underset{j=1}{\overset{m}{}}\xi _jq_j,q^2=\underset{j=1}{\overset{m}{}}q_j^2.$$
If a Gaussian measure $`\mu `$ is nondegenerate (so the measure of any open set is positive), then
$$d\mu (q)=\frac{e^{\frac{1}{2}(B^1q,q)}dq}{\sqrt{(2\pi )^mdetB}},$$
where $`B`$ is a positive operator (we consider everywhere only Gaussian measures with zero mean values). If $`\mu S_G^\alpha (Q)`$ and nondegenerate, then
$$\sigma ^2(\mu )=\frac{1}{\sqrt{(2\pi )^mdetB}}_{𝐑^m}q^2e^{\frac{1}{2}(B^1q,q)}𝑑q=\alpha +O(\alpha ^2).$$
In the general case the easiest way to define a Gaussian measure is to use its Fourier transform:
$$\stackrel{~}{\mu }(\xi )=_{𝐑^m}e^{i(\xi ,q)}𝑑\mu (q)=e^{\frac{1}{2}(Bq,q)},$$
where $`B=\mathrm{cov}\mu `$ is the covariance operator:
$$(B\xi _1,\xi _2)=_{𝐑^m}(\xi _1,q)(\xi _2,q)𝑑\mu (q).$$
We remark that by definition a covariance operator is positively defined and symmetric.
Lemma 5.1. Let $`\mu `$ be a Gaussian measure with the zero mean value and let $`A`$ be a symmetric operator. Then
$$_{𝐑^m}(Aq,q)𝑑\mu (q)=\mathrm{Tr}\mathrm{BA},$$
(29)
where $`B=\mathrm{cov}\mu .`$
To prove this lemma we should just expand the quadratic form $`(Aq,q)`$ with respect to an orthonormal basis.
Corollary 5.1. We have
$$\sigma ^2(\mu )=_{𝐑^m}q^2𝑑\mu (q)=\mathrm{Tr}\mathrm{B}.$$
(30)
Thus, for $`\mu S_G^\alpha (Q),`$
$$\mathrm{Tr}\mathrm{cov}\mu =\alpha +\mathrm{O}(\alpha ^2).$$
We now define a class of physical variables – $`𝒱(Q):`$ a) $`fC^{\mathrm{}}(𝐑^m);`$ b) $`f(0)=0;`$ c) $`f^{(4)}(q)c_fe^{r_fq},c_f,r_f0.`$
For a function $`f:𝐑^m𝐑,`$ its $`n`$th derivative is a (symmetric) $`n`$-linear functional, $`f^{(n)}(q):𝐑^m\times \mathrm{}\times 𝐑^m𝐑.`$ The norm of this functional is given by
$$f^{(n)}(q)=\underset{h_j=1}{sup}|f^{(n)}(q)(h_1,\mathrm{},h_n)|.$$
The norm can be estimated by partial derivatives:
$$f^{(n)}(q)\underset{\alpha _1+\mathrm{}+\alpha _n=n}{\mathrm{max}}\left|\frac{^nf(q)}{q_1^{\alpha _1}\mathrm{}q_n^{\alpha _n}}\right|.$$
It is easy to generalize Lemma 4.1 to the multidimensional case.
Thus we have defined the following classical statistical model: $`N_{\mathrm{class}}=(S_G^\alpha (Q),𝒱(Q)).`$
Lemma 5.2. Let $`f𝒱(Q)`$ and let $`\mu S_G^\alpha (Q).`$ Then
$$<f>_\mu (\alpha )_{𝐑^m}f(q)𝑑\mu (q)=\frac{\alpha }{2}\mathrm{Tr}\rho \mathrm{f}^{\prime \prime }(0)+\mathrm{O}(\alpha ^2),$$
(31)
where $`\rho `$ is a density operator; in fact, $`\rho =\mathrm{cov}\mu /\alpha .`$
Proof. By using the scaling of the state variable (22) and by expanding $`f(\sigma (\mu )x)`$ on the basis of the fourth order Taylor formula with the integral remainder we obtain:
$$<f>_\mu (\alpha )=\frac{\sigma ^2(\mu )}{2}\mathrm{Tr}\rho \mathrm{f}^{\prime \prime }(0)$$
(32)
$$+\frac{\sigma ^4(\mu )}{4!}_{𝐑^m}\left(_0^1(1\theta )^3f^{(4)}(\sigma (\mu )x\theta )(q,q,q,q)𝑑\theta \right)𝑑\mu _{\mathrm{scal}}(x),$$
where $`\mu _{\mathrm{scal}}`$ is a normalized Gaussian measure – the image of $`\mu `$ under the scaling (22). We now estimate the remainder:
$$|R(f,\mu )|\frac{C\sigma ^4(\mu )}{4!}_{𝐑^m}x^4\left(_0^1(1\theta )^3e^{r\sigma (\mu )x\theta }𝑑\theta \right)𝑑\mu _{\mathrm{scal}}(x).$$
Thus
$$|R(f,\mu )|C^{}\sigma ^4(\mu )_{𝐑^m}x^4e^{rx}𝑑\mu _{\mathrm{scal}}(x).$$
We have that $`R(f,\mu )=O(\alpha ^2),\alpha 0.`$
Corollary 5.2. Let $`f𝒱(Q)`$ and let $`\mu S_G^\alpha (Q)`$ be nondegenerate. Then
$$\frac{1}{\sqrt{(2\pi )^mdetB}}_{𝐑^m}f(q)e^{\frac{1}{2}(B^1q,q)}𝑑q=\frac{\alpha }{2}\mathrm{Tr}\rho \mathrm{f}^{\prime \prime }(0)+\mathrm{O}(\alpha ^2).$$
(33)
where $`\rho =B/\alpha .`$
As in the one-dimensional case, we introduce the relative average:
$$f_\mu \frac{<f>_\mu }{\alpha }=\frac{_{𝐑^m}f(q)𝑑\mu (q)}{_{𝐑^m}q^2𝑑\mu (q)}+O(\alpha ).$$
In the case of a nondegenerate Gaussian measure we have:
$$f_\mu =\frac{_{𝐑^m}f(q)e^{\frac{1}{2}(B^1q,q)}𝑑q}{_{𝐑^m}q^2e^{\frac{1}{2}(B^1q,q)}𝑑q}+O(\alpha ).$$
Corollary 5.3. Let $`f𝒱(Q)`$ and let $`\mu S_G^\alpha (Q).`$ Then
$$f_\mu =\frac{1}{2}\mathrm{Tr}\rho \mathrm{f}^{\prime \prime }(0)+\mathrm{O}(\alpha ).$$
(34)
Thus if one neglects by terms of the magnitude $`\alpha ,`$ it is possible to use the following approximative calculus of averages:
$$f_\mu ^{\mathrm{approx}}=\frac{1}{2}\mathrm{Tr}\rho \mathrm{A},$$
(35)
where $`A=f^{\prime \prime }(0)`$ and $`\rho =\mathrm{cov}\mu _{\mathrm{scal}}.`$ This is nothing else than the von Neumann trace formula for quantum averages, see (13). To proceed more formally, we consider maps:
$$T:S_G^\alpha (Q)𝒟^{(r)}(m),\rho =T(\mu )=\mathrm{cov}\mu _{\mathrm{scal}};$$
(36)
$$T:𝒱(Q))_s^{(r)}(m)),A=T(f)=f^{\prime \prime }(0)$$
(37)
(we recall that Hessian is always a symmetric matrix).
Theorem 5.1. The maps (36), (37) project the classical statistical model $`N_{\mathrm{class}}=(S_G^\alpha (Q),𝒱(Q))`$ onto the quantum model $`N_{\mathrm{quant}}^{(r)}=(𝒟^{(r)}(m),_s^{(r)}(m))`$ in such a way that classical and quantum averages are coupled by the asymptotic equality:
$$f_\mu =\frac{1}{2}<T(f)>_{T(\mu )}+O(\alpha ).$$
(38)
## 6 Degenerate Gaussian measures and “pure states”
Consider a Gaussian measure $`\mu S_G^\alpha (Q)`$ which is concentrated on a linear subspace $`Q_0`$ of $`Q.`$ So it is nondegenerate on $`Q_0.`$ Here $`P:QQ_0`$ is the orthogonal projector onto $`Q_0.`$ Denote by $`B_0`$ the covariance matrix of the restriction of $`\mu `$ onto $`Q_0.`$ Then $`B_0>0`$ and $`B=\mathrm{cov}\mu =\mathrm{PB}_0\mathrm{P}.`$ Let us now make the scaling (22). Then $`\rho =\mathrm{cov}\mu _{\mathrm{scal}}=\mathrm{P}\rho _0\mathrm{P},`$ where $`\rho _0=\mathrm{cov}\mu _{\mathrm{scal}}|_{\mathrm{Q}_0}.`$
Thus for any symmetric matrix $`A`$ we have $`\mathrm{Tr}\rho \mathrm{A}=\mathrm{Tr}\rho _0(\mathrm{PAP}).`$ Suppose now that $`\rho _0=\frac{I}{\mathrm{dim}\mathrm{Q}_0},`$ where $`I:Q_0Q_0`$ is the unit operator. Then $`\mathrm{Tr}\rho \mathrm{A}=\frac{1}{\mathrm{dim}\mathrm{Q}_0}\mathrm{Tr}\mathrm{PAP}.`$ We are especially interested in measures concentrated on one dimensional subspaces, $`Q_0Q_\mathrm{\Psi }=\{q=c\mathrm{\Psi },c𝐑\},`$ where $`\mathrm{\Psi }`$ has the norm one. Here $`PP_\mathrm{\Psi }`$ is the one dimensional projector $`P_\mathrm{\Psi }:QQ_\mathrm{\Psi },P_\mathrm{\Psi }\varphi =(\varphi ,\mathrm{\Psi })\mathrm{\Psi },`$ and hence
$$\mathrm{Tr}\rho \mathrm{A}=\mathrm{Tr}\mathrm{P}_\mathrm{\Psi }\mathrm{AP}_\mathrm{\Psi }=(\mathrm{A}\mathrm{\Psi },\mathrm{\Psi }).$$
(39)
Denote a probability $`\mu S_G^\alpha (Q)`$ which is concentrated on the one dimensional subspace $`Q_\mathrm{\Psi }`$ by the symbol $`\mu _\mathrm{\Psi }.`$ We obtained the following simple result:
Proposition 6.1. For any $`\mu _\mathrm{\Psi },`$ we have
$$f_{\mu _\mathrm{\Psi }}=\frac{1}{2}(f^{\prime \prime }(0)\mathrm{\Psi },\mathrm{\Psi })+O(\alpha ).$$
(40)
Thus approximately we have:
$$f_{\mu _\mathrm{\Psi }}=\frac{1}{2}(A\mathrm{\Psi },\mathrm{\Psi }),A=f^{\prime \prime }(0).$$
(41)
But the right-hand side of this equality is nothing else than the well known quantum formula for the average of the quantum observable $`A`$ with respect to the pure state $`\mathrm{\Psi },`$ . In our approach this quantum formula arose as the approximation of the classical average with respect to a Gaussian ensemble. The only distinguishing feature of such an ensemble is that all systems have states proportional to the vector $`\mathrm{\Psi }`$ (with the probability one). Of course, one can consider the projective space and then all those systems will have the same coordinate. However, real coordinates of systems are different.
Conclusion. Quantum averages with respect so called pure states can be easily reproduced as approximations of ordinary ensemble averages with respect to one dimensional Gaussian distributions.
## 7 Prequantum phase space – the two dimensional case
In previous sections we considered the prequantum toy model in that the phase space structure was not taken into account. The corresponding quantum model was over the reals, see also , . On the other hand, physical reality is described by the classical phase space mechanics and the complex quantum mechanics. We shall see that it is possible to create a prequantum phase space model reproducing the complex quantum mechanics. The crucial point is that classical variables and statistical states – functions and measures on phase space – should be invariant with respect to a special group of transformations of phase space.
This fundamental prequantum group is very simple – the special orthogonal group $`SO(2),`$ the group of rotations of phase space.
States of systems are now represented by points $`\psi =(q,p)\mathrm{\Omega }=Q\times P,`$ where $`Q=P=𝐑.`$ Here the $`q`$ is the position and the $`p`$ is momentum, so $`\mathrm{\Omega }`$ denotes phase space. Statistical states are represented by Gaussian $`SO(2)`$-invariant measures having zero mean value and dispersion
$$\sigma ^2(\mu )=2\alpha +O(\alpha ^2);$$
(42)
physical variables are by $`SO(2)`$-invariant maps, $`f:\mathrm{\Omega }𝐑,`$ which satisfy conditions a), b), c) specifying variables in the real case. Denote these classes of measures and functions, respectively, $`S_G^\alpha (\mathrm{\Omega }|SO(2))`$ and $`𝒱(\mathrm{\Omega }|SO(2)).`$
The appearance of the factor 2 has the following motivation: there are two contributions into fluctuations – fluctuations of positions and momenta. We shall see that they are equally distributed. Therefore it is natural to consider as a small parameter of the model the dispersion of e.g. the $`q`$-fluctuations (which equals to the dispersion of the $`p`$-fluctuations).
We consider the classical model $`N_{\mathrm{class}}=(S_G^\alpha (\mathrm{\Omega }|SO(2)),𝒱(\mathrm{\Omega }|SO(2)).`$ As in the real case, we can obtain the asymptotic expansion of the classical averages, see (34). However, in quantum mechanics we consider the complex structure. We would like to recover it in our classical model. To do this, we shall study in more detail properties of classical probabilities and variables.
A measure $`\mu `$ is invariant if for any $`uSO(2):`$
$$_{𝐑^2}f(uq)𝑑\mu (q)=_{𝐑^2}f(q)𝑑\mu (q).$$
(43)
For a Gaussian measure $`\mu `$ with the covariance matrix $`B,`$ this is equivalent to the condition:
$$[u,B]=0,uSO(2).$$
(44)
Let $`f`$ be a two times differentiable invariant map, so $`f(u\psi )=f(\psi ),`$ for any $`uSO(2).`$ By representing
$$u=u_\theta =\left(\begin{array}{cc}\mathrm{cos}\theta \hfill & \mathrm{sin}\theta \hfill \\ \mathrm{sin}\theta \hfill & \mathrm{cos}\theta \hfill \end{array}\right),$$
(45)
we have that
$$f(\mathrm{cos}\theta q\mathrm{sin}\theta p,\mathrm{sin}\theta q+\mathrm{cos}\theta p)=f(q,p).$$
(46)
This is a rather strong constraint determining a very special class of maps. In particular, we obtain: $`u^{}f(u\psi )=f(\psi )`$ and $`u^{}f^{\prime \prime }(u\psi )u=f^{\prime \prime }(\psi ).`$ Hence $`u^{}f(0)=f(0)`$ for any rotation, and thus
$$f(0)=0$$
(47)
and
$$[f^{\prime \prime }(0),u]=0,uSO(2).$$
(48)
It is convenient to introduce the commutator of the set $`SO(2)`$ in the algebra of all two by two matrices $`M^{(r)}(2):`$
$$SO^{}(2)=\{AM^{(r)}(2):[A,u]=0,uSO(2)\}$$
We remark that a generator of $`SO(2)`$ can be chosen as the symplectic operator:
$$J=\left(\begin{array}{cc}0\hfill & 1\hfill \\ 1\hfill & 0\hfill \end{array}\right).$$
Therefore the commutator of $`SO^{}(2)`$ coincides with the commutator of $`J:\{J\}^{}=\{AM^{(r)}(2):[A,J]=0\}.`$
Proposition 7.1. Let $`\mu S_G^\alpha (\mathrm{\Omega }|SO(2))`$and let $`f𝒱(\mathrm{\Omega }|SO(2)).`$ Then $`B=\mathrm{cov}\mu `$ and $`A=f^{\prime \prime }(0)`$ belong to $`SO^{}(2).`$
Lemma 7.1. A matrix $`A`$ belongs to the commutator $`SO^{}(2)`$ iff
$$A=\left(\begin{array}{cc}R\hfill & S\hfill \\ S\hfill & R\hfill \end{array}\right).$$
(49)
If $`A`$ is also symmetric, then it is diagonal: $`A=\left(\begin{array}{cc}R\hfill & 0\hfill \\ 0\hfill & R\hfill \end{array}\right).`$ In particular, its trace is given by
$$\mathrm{TrA}=2\mathrm{R}$$
(50)
Thus if $`\mu S_G^\alpha (\mathrm{\Omega }|SO(2)),`$ then its covariance matrix is diagonal $`B=\left(\begin{array}{cc}b\hfill & 0\hfill \\ 0\hfill & b\hfill \end{array}\right),`$ where $`2b=\alpha +O(\alpha ^2).`$ Fluctuations of the coordinate $`q`$ and the momentum $`p`$ are independent and equally distributed:
$$d\mu (q)=\frac{1}{2\pi b}\mathrm{exp}\{\frac{q^2+p^2}{2b}\}dq.$$
Denote the marginal distributions of $`\mu `$ by the symbols $`\mu _q`$ and $`\mu _p,`$ respectively. Then
$$\sigma ^2(\mu _q)=\frac{1}{\sqrt{2\pi b}}_{\mathrm{}}^+\mathrm{}q^2\mathrm{exp}\{\frac{q^2}{2b}\}𝑑q=\sigma ^2(\mu _p)=\frac{1}{\sqrt{2\pi b}}_{\mathrm{}}^+\mathrm{}p^2\mathrm{exp}\{\frac{p^2}{2b}\}𝑑p.$$
Hence
$$\sigma ^2(\mu _q)=\sigma ^2(\mu _p)=\frac{1}{2}\sigma ^2(\mu )=\alpha +O(\alpha ^2).$$
Proposition 7.2. Let $`f𝒱(\mathrm{\Omega }|SO(2)).`$ Then all its odd derivatives at the point $`q_0=0`$ are equal to zero.
Proof. For any vector $`\varphi \mathrm{\Omega },`$ we have $`f^{(2n+1)}(u\psi )(u\varphi ,\mathrm{},u\varphi )`$
$`=f^{(2n+1)}(\psi )(\varphi ,\mathrm{},\varphi ),`$ for any rotation $`u.`$ We choose $`u=J.`$ Then:
$$(1)^{2n+1}f^{(2n+1)}(0)(\varphi ,\mathrm{},\varphi )=f^{(2n+1)}(0)(J^2\varphi ,\mathrm{},J^2\varphi )$$
$$=f^{(2n+1)}(0)(J\varphi ,\mathrm{},J\varphi )=f^{(2n+1)}(0)(\varphi ,\mathrm{},\varphi ).$$
Thus $`f^{(2n+1)}(0)(\varphi ,\mathrm{},\varphi )=0`$ for any vector $`\varphi \mathrm{\Omega }.`$
For a function $`f𝒱(\mathrm{\Omega }|SO(2)),`$ its Hessian has the form $`f^{\prime \prime }(0)=\left(\begin{array}{cc}R\hfill & 0\hfill \\ 0\hfill & R\hfill \end{array}\right),`$ where $`R𝐑,`$ and hence:
$$f(q,p)=\frac{R(q^2+p^2)}{2}+O(\alpha ^2).$$
We remark that, in spite of the coincidence of commutators, the $`SO(2)`$-invariance is not equivalent to the $`J`$-invariance (the later was used as the basis of the theory in ).
Example 7.1. Let $`f(q,p)=q^3pqp^3=qp(q^2p^2).`$ Then $`f(J\psi )=f(\psi ).`$ But take $`\theta =\pi /4.`$ Here $`u=\frac{1}{2}\left(\begin{array}{cc}1\hfill & 1\hfill \\ 1\hfill & 1\hfill \end{array}\right).`$ Hence $`u\left(\begin{array}{c}q\hfill \\ p\hfill \end{array}\right)=\left(\begin{array}{c}(qp)/2\hfill \\ (q+p)/2\hfill \end{array}\right).`$ Thus $`f(u\psi )=(qp)(q+p)(qpqp)(qp+q+p)/16=qp((q^2p^2)/4.`$
We are now completely ready to recover the complex structure of quantum mechanics. By Lemma 7.1 any matrix belonging $`SO^{}(2)`$ can be represented in the form: $`A=RI+S(J).`$ By mapping $`I`$ into 1 and $`(J)`$ into $`i`$ we obtain a map of the commutator $`SO^{}(2)`$ onto the set of complex numbers $`𝐂:`$
$$j:SO^{}(2)𝐂,z=j(A)=R+iS.$$
(51)
This is the isomorphism of two fields.
In particular, a symmetric matrix $`A=\left(\begin{array}{cc}R\hfill & 0\hfill \\ 0\hfill & R\hfill \end{array}\right)`$ is represented by the real number $`j(A)=R.`$ This is the operator of multiplication by $`R.`$ The trace of this operator in the one dimensional complex space $`𝐂`$ (with the scalar product, $`(z,w)=z\overline{w})`$ equals $`R.`$ By (50) we obtain that
$$\mathrm{Tr}\mathrm{A}=2\mathrm{T}\mathrm{r}\mathrm{j}(\mathrm{A}),$$
(52)
where at the left-hand side we have the real trace and at the right-hand side – the complex trace. Now we can write the basic asymptotic equality for averages in the complex form. In the funny way the Taylor factor $`\frac{1}{2}`$ disappears through the transition form the real to complex structure, see (52).
Lemma 7.2. Let $`f𝒱(\mathrm{\Omega }|SO(2))`$ and let $`\mu S_G^\alpha (\mathrm{\Omega }|SO(2)).`$ Then
$$<f>_\mu (\alpha )_{𝐑^2}f(q,p)𝑑\mu (q,p)=\alpha j(f^{\prime \prime }(0))+O(\alpha ^2).$$
(53)
Proof. We make the scaling of the state variable:
$$\psi =\frac{\sigma (\mu )}{\sqrt{2}}\mathrm{\Psi }$$
(54)
Then the image of $`\mu `$ is again a Gaussian measure, say $`\mu _{\mathrm{scal}},`$ having the dispersion $`\sigma ^2(\mu _{\mathrm{scal}})=2.`$ Set $`D=\mathrm{cov}\mu _{\mathrm{scal}}.`$ In the two dimensional case $`D=\left(\begin{array}{cc}1\hfill & 0\hfill \\ 0\hfill & 1\hfill \end{array}\right)`$ and $`\mathrm{Tr}\mathrm{D}=2.`$ We now have:
$$<f>_\mu (\alpha )=\frac{\sigma ^2(\mu )}{4}\mathrm{Tr}\mathrm{D}\mathrm{f}^{\prime \prime }(0)+\mathrm{O}(\alpha ^2).$$
(55)
Thus
$$<f>_\mu (\alpha )=\frac{\sigma ^2(\mu )}{2}j(D)j(f^{\prime \prime }(0))+O(\alpha ^2).$$
(56)
Finally, we note that in the two dimensional case: $`j(D)=1.`$ Thus we obtain:
$$<f>_\mu (\alpha )=\frac{\sigma ^2(\mu )}{2}j(f^{\prime \prime }(0))+O(\alpha ^2),$$
(57)
and hence (53).
We recall that in the one dimensional quantum mechanics there is just one “density matrix”, namely, $`\rho =1𝐑.`$
It is convenient to consider the renormalization of averages by the main term in the intensities of fluctuations of the coordinate and momenta: $`f_\mu =\frac{<f>_\mu }{\alpha }.`$ Then we get:
$$f_\mu (\alpha )=j(f^{\prime \prime }(0))+O(\alpha ).$$
(58)
## 8 Prequantum phase space – multidimensional case
States of systems are now represented by points $`\psi =(q,p)\mathrm{\Omega }=Q\times P,`$ where $`Q=P=𝐑^m.`$ Here the $`q=(q_1,\mathrm{},q_n)`$ is the position and the $`p=(p_1,\mathrm{},p_n)`$ is momentum, so $`\mathrm{\Omega }`$ denotes phase space. Let us consider the canonical representation of the group $`SO(2)`$ in the phase space $`\mathrm{\Omega }=Q\times P:`$
$$u=u_\theta =\left(\begin{array}{cc}\mathrm{cos}\theta I\hfill & \mathrm{sin}\theta I\hfill \\ \mathrm{sin}\theta I\hfill & \mathrm{cos}\theta I\hfill \end{array}\right),$$
(59)
where $`I`$ is the unit matrix from $`M^{(r)}(m).`$ The corresponding group of $`𝐑`$-linear operators (or $`2m\times 2m`$ matrices) we denote by the symbol $`SO_m(2).`$
The classical model $`N_{\mathrm{class}}=(S_G^\alpha (\mathrm{\Omega }|SO_m(2)),𝒱(\mathrm{\Omega }|SO_m(2)))`$ in defined in the same way as in the two dimensional case. A Gaussian measure is invariant iff its covariance operator belongs to the commutator $`SO_m^{}(2)=\{AM^{(r)}(2m):[A,u]=0,uSO_m(2)\}.`$ If a smooth function $`f`$ is invariant then all its odd derivatives equal to zero and the second derivative belong to the $`SO_m^{}(2).`$ A matrix $`ASO_m^{}(2)`$ if it has the form (49), where $`R,SM^{(r)}(m).`$ In contrast to the two dimensional case a symmetric matrix from $`SO_m^{}(2)`$ can be nondiagonal. It has the form (49), where $`R^{}=R`$ and $`S^{}=S.`$
There is a natural map (generalizing the map $`j:SO^{}(2)𝐂)`$ of the commutator $`SO_m^{}(2)`$ onto the set of complex matrices $`M^{(c)}(m):`$
$$j:SO_m^{}(2)M^{(c)}(m),z=j(A)=R+iS.$$
(60)
This is the isomorphism of two rings. Symmetric matrices are mapped onto symmetric matrices. Let us denote real and complex conjugations by $``$ and $`,`$ respectively. We have $`(R+iS)^{}=R^{}iS^{}=R+iS.`$ We also remark that for a symmetric complex matrix:
$$\mathrm{Tr}\mathrm{j}(\mathrm{A})=\mathrm{Tr}(\mathrm{R}+\mathrm{iS})=\mathrm{Tr}\mathrm{R}=\frac{1}{2}\mathrm{Tr}\mathrm{A}.$$
(61)
Lemma 8.1. Let $`f𝒱(\mathrm{\Omega }|SO_m(2))`$ and let $`\mu S_G^\alpha (\mathrm{\Omega }|SO_m(2)).`$ Then
$$<f>_\mu (\alpha )=\alpha \mathrm{Tr}\rho \mathrm{j}(\mathrm{f}^{\prime \prime }(0))+\mathrm{O}(\alpha ^2),$$
(62)
where $`\rho 𝒟^{(c)}(m).`$
Proof. We make the scaling (54) and get the $`\mu _{\mathrm{scal}}`$ with $`D=\mathrm{cov}\mu _{\mathrm{scal}},`$ and $`\mathrm{Tr}\mathrm{D}=2.`$ We set $`\rho =j(D),`$ here $`\mathrm{Tr}\rho =(\mathrm{Tr}\mathrm{D})/2=1`$ and $`\rho 𝒟^{(c)}(m).`$ Finally
$$<f>_\mu (\alpha )=\frac{\sigma ^2(\mu )}{2}\mathrm{Tr}\mathrm{j}(\mathrm{D})\mathrm{j}(\mathrm{f}^{\prime \prime }(0))+\mathrm{O}(\alpha ^2)$$
(63)
implies (62).
We now modify the classical$``$ quantum projections, (36), (37), to make them consistent with the complex structure:
$$T:S_G^\alpha (\mathrm{\Omega }|SO_m(2))𝒟^{(c)}(m),\rho =T(\mu )=j(\mathrm{cov}\mu _{\mathrm{scal}});$$
(64)
$$T:𝒱(\mathrm{\Omega }|SO_m(2))_s^{(c)}(m),A=T(f)=j(f^{\prime \prime })(0)$$
(65)
Theorem 8.1. The maps (64), (65) project the classical statistical model $`N_{\mathrm{class}}=(S_G^\alpha (\mathrm{\Omega }|SO_m(2)),𝒱(\mathrm{\Omega }|SO_m(2)))`$ onto the quantum model $`N_{\mathrm{quant}}^{(c)}=(𝒟^{(c)}(m),_s^{(c)}(m))`$ in such a way that classical and quantum averages are coupled by the asymptotic equality:
$$f_\mu =<T(f)>_{T(\mu )}+O(\alpha ).$$
(66)
## 9 Prequantum phase space
States of systems are now represented by points $`\psi =(q,p)\mathrm{\Omega }=Q\times P,`$ where $`Q=P=H`$ and $`H`$ is a real (separable Hilbert space) with the scalar product $`(,)`$ and the corresponding norm $`.`$ Here the $`qH`$ is the position and the $`pH`$ is momentum, so $`\mathrm{\Omega }`$ denotes phase space. The real Hilbert space structure on $`\mathrm{\Omega }`$ is given by the scalar product
$$(\psi _1,\psi _2)=(q_1,q_2)+(p_1,p_2).$$
(67)
In physics $`H=L_2(𝐑^3)`$ is the space of square integrable functions. Thus both position and momentum are functions of $`x𝐑^3.`$ These are simply classical fields. A point of such a phase space is a classical vector field $`\psi (x)=(q(x),p(x)).`$
Let us consider the canonical representation of the group $`SO(2)`$ in the phase space $`\mathrm{\Omega }=Q\times P,`$ see (59), where $`I:HH`$ is the unit operator. The corresponding group of continuous $`𝐑`$-linear operators we denote by the symbol $`SO_H(2).`$
The classical model
$$N_{\mathrm{class}}=(S_G^\alpha (\mathrm{\Omega }|SO_H(2)),𝒱(\mathrm{\Omega }|SO_H(2)))$$
in defined in the same way as in the finite-dimensional case. We just recall a few basic notions from theory of differentiable functions and Gaussian measures on infinite-dimensional spaces.
Let $`\mu `$ be a $`\sigma `$-additive Gaussian measure on the $`\sigma `$-field $`F`$ of Borel subsets of $`\mathrm{\Omega },`$ see . This measure is determined by its covariance operator $`B:\mathrm{\Omega }\mathrm{\Omega }`$ and mean value $`mm_\mu \mathrm{\Omega }.`$ For example, $`B`$ and $`m`$ determines the Fourier transform of $`\rho :`$
$$\stackrel{~}{\rho }(y)=_\mathrm{\Omega }e^{i(y,\psi )}𝑑\mu (\psi )=e^{\frac{1}{2}(By,y)+i(m,y)},y\mathrm{\Omega }.$$
In what follows we restrict our considerations to Gaussian measures with zero mean value $`m=0,`$ where $`(m,y)=_\mathrm{\Omega }(y,\psi )𝑑\mu (\psi )=0`$ for any $`y\mathrm{\Omega }.`$ We recall that the covariance operator $`B\mathrm{cov}\mu `$ is defined by
$$(By_1,y_2)=_\mathrm{\Omega }(y_1,\psi )(y_2,\psi )𝑑\mu (\psi ),y_1,y_2\mathrm{\Omega },$$
and has the following properties: a). $`B0,`$ i.e., ($`By,y)0,y\mathrm{\Omega };`$ b). $`B`$ is a bounded self-adjoint operator, $`B_s(\mathrm{\Omega });`$ c). $`B`$ is a trace-class operator and moreover
$$\mathrm{Tr}B=_\mathrm{\Omega }\psi ^2𝑑\mu (\psi ).$$
This is dispersion $`\sigma ^2(\mu )`$ of the probability $`\mu .`$ Thus $`\sigma ^2(\mu )=\mathrm{Tr}B.`$
We remark that the list of properties of the covariance operator of a Gaussian measure differs from the list of properties of a von Neumann density operator only by one condition: $`\mathrm{Tr}\rho =1,`$ for a density operator $`\rho .`$
We can easily find the Gaussian integral of a quadratic form (by using expansion with respect to an orthonormal basis and using our previous results on the finite-dimensional Gaussian integrals): $`_\mathrm{\Omega }(A\psi ,\psi )𝑑\rho (\psi )=\mathrm{Tr}BA,`$ where $`A_s(\mathrm{\Omega }).`$
The differential calculus for maps $`f:\mathrm{\Omega }𝐑`$ does not differ so much from the differential calculus in the finite dimensional case, $`f:𝐑^n𝐑.`$ Instead of the norm on $`𝐑^n,`$ one should use the norm on $`\mathrm{\Omega }.`$ We consider so called Frechet differentiability. Here a function $`f`$ is differentiable if it can be represented as
$$f(\psi _0+\mathrm{\Delta }\psi )=f(\psi _0)+f^{}(\psi _0)(\mathrm{\Delta }\psi )+o(\mathrm{\Delta }\psi ),$$
where $`lim_{\mathrm{\Delta }\psi 0}\frac{o(\mathrm{\Delta }\psi )}{\mathrm{\Delta }\psi }=0.`$ Here at each point $`\psi `$ the derivative $`f^{}(\psi )`$ is a continuous linear functional on $`\mathrm{\Omega };`$ so it can be identified with the element $`f^{}(\psi )\mathrm{\Omega }.`$ Then we can define the second derivative as the derivative of the map $`\psi f^{}(\psi )`$ and so on. A map $`f`$ is differentiable $`n`$-times iff:
$$f(\psi _0+\mathrm{\Delta }\psi )=f(\psi _0)+f^{}(\psi _0)(\mathrm{\Delta }\psi )+\frac{1}{2}f^{\prime \prime }(\psi _0)(\mathrm{\Delta }\psi ,\mathrm{\Delta }\psi )+\mathrm{}$$
$$+\frac{1}{n!}f^{(n)}(\psi _0)(\mathrm{\Delta }\psi ,\mathrm{},\mathrm{\Delta }\psi )+o_n(\mathrm{\Delta }\psi ),$$
where $`f^{(n)}(\psi _0)`$ is a symmetric continuous $`n`$-linear form on $`\mathrm{\Omega }`$ and
$$\underset{\mathrm{\Delta }\psi 0}{lim}\frac{o_n(\mathrm{\Delta }\psi )}{\mathrm{\Delta }\psi ^n}=0.$$
For us it is important that the second derivative $`f^{\prime \prime }(\psi _0)`$ can be represented by a self-adjoint operator $`f^{\prime \prime }(\psi _0)(u,v)=(f^{\prime \prime }(\psi _0)u,v),u,v\mathrm{\Omega }.`$ We remark that for $`\psi _0=0`$ we have:
$$f(\psi )=f(0)+f^{}(0)(\psi )+\frac{1}{2}f^{\prime \prime }(0)(\psi ,\psi )+\mathrm{}+\frac{1}{n!}f^{(n)}(0)(\psi ,\mathrm{},\psi )+o_n(\psi ).$$
As in the finite-dimensional case the reminder can be represented in the integral form.
A Gaussian measure is invariant iff its covariance operator belongs to the commutator
$$SO_H^{}(2)=\{A(\mathrm{\Omega }):[A,u]=0,uSO_H(2)\}.$$
If a smooth function $`f:\mathrm{\Omega }𝐑`$ is $`SO_H^{}(2)`$-invariant then all its odd derivatives equal to zero and the second derivative belong to the $`SO_H^{}(2).`$ An operator $`ASO_H^{}(2)`$ if it has the form (49), where $`R,S(H).`$ And $`ASO_H^{}(2)_s(\mathrm{\Omega })`$ if it has the form (49), where $`R^{}=R`$ and $`S^{}=S.`$
Let us now consider the complexification of the Hilbert space $`H:`$ $`H_c=HiH.`$ We denote the algebra of bounded $`𝐂`$-linear operators $`A:H_cH_c`$ by the symbol $`(H_c).`$ The set of self-adjoint operators (with respect to the complex scalar product) we denote by the symbol $`_s(H_c).`$
There is a natural map (generalizing the map $`j:SO_m^{}(2)𝐂^m)`$ of the commutator $`SO_H^{}(2)`$ onto the $`(H_c):`$
$$j:SO_H^{}(2)(H_c),z=j(A)=R+iS.$$
(68)
This is the isomorphism of two rings. Self-adjoint operators (with respect to the real scalar product) are mapped onto self-adjoint operators (with respect to the complex scalar product). We also remark that for a self-adjoint trace class operator $`ASO_H^{}(2)`$ the equality (61) coupling real and comlex traces holds. In the same way as in the finite-dimensional case we prove:
Lemma 9.1. Let $`f𝒱(\mathrm{\Omega }|SO_H(2))`$ and let $`\mu S_G^\alpha (\mathrm{\Omega }|SO_H(2)).`$ Then the asymptotic equality (62), where $`\rho 𝒟(H_c),`$ holds.
We now consider the infinite-dimensional generalization of the classical$``$ quantum projections, (64), (65)
$$T:S_G^\alpha (\mathrm{\Omega }|SO_H(2))𝒟(H_c),\rho =T(\mu )=j(\mathrm{cov}\mu _{\mathrm{scal}});$$
(69)
$$T:𝒱(\mathrm{\Omega }|SO_H(2))_s(H_c),A=T(f)=j(f^{\prime \prime })(0)$$
(70)
Lemma 9.1 implies:
Theorem 9.1. The maps (69), (70) project the classical statistical model $`N_{\mathrm{class}}=(S_G^\alpha (\mathrm{\Omega }|SO_H(2)),𝒱(\mathrm{\Omega }|SO_H(2)))`$ onto the quantum model $`N_{\mathrm{quant}}(H_c)=(𝒟(H_c),_s(H_c))`$ in such a way that classical and quantum averages are coupled by the asymptotic equality (66).
We remark that the idea that the quantum averages can be coupled to integration with respect to the $`\psi `$-function was discussed in a number of papers, see, e.g., and (so called GAP-measures) as well as extended literature in the last paper. The main distinguishing feature of our approach is elaboration of technique of asymptotic expansion with respect to a small parameter, namely the dispersion of prequantum fluctuations. Comparing with we also mention that we (as well as Bach ) consider the linear space integration and not integration over the unit sphere. This is not simply a technical deviation, but it implies a totally new viewpoint to quantum pure states, see the next section.
## 10 Gaussian measures corresponding to pure quantum states
We now generalize considerations of section 6 to the complex infinite-dimensional case. Let $`\mathrm{\Psi }`$ be a pure quantum state: $`\mathrm{\Psi }H_c,\mathrm{\Psi }=1.`$ We define a Gaussian measure $`\mu _\mathrm{\Psi }`$ which is concentrated on the one dimensional complex space (so the real plane) $`\mathrm{\Pi }_\mathrm{\Psi }=\{\psi H_C:\psi =c\mathrm{\Psi },c𝐂\}:`$ the average of the $`\mu _\mathrm{\Psi }`$ is equal to zero and the complexification of its covariance operator:
$$j(\rho _{\mu _\mathrm{\Psi }})=\alpha \mathrm{\Psi }\mathrm{\Psi }.$$
The following facts about $`\mu _\mathrm{\Psi }`$ can be obtained through direct computations and Theorem 9.1:
Proposition 10.1. For any pure quantum state $`\mathrm{\Psi },`$ we have:
a). $`\mu _\mathrm{\Psi }S_G^\alpha (\mathrm{\Omega }|SO_H(2));`$
b). $`T(\mu _\mathrm{\Psi })=\mathrm{\Psi }\mathrm{\Psi };`$
c). $`f_\mu =<j(f^{\prime \prime }(0))\mathrm{\Psi },\mathrm{\Psi }>+O(\alpha ),f𝒱(\mathrm{\Omega }|SO_H(2)).`$
## 11 Hamilton-Schrödinger dynamics
States of systems with the infinite number of degrees of freedom - classical fields – are represented by points $`\psi =(q,p)\mathrm{\Omega };`$ evolution of a state is described by the Hamiltonian equations. We consider a quadratic Hamilton function: $`(q,p)=\frac{1}{2}(𝐇\psi ,\psi ),`$ where $`𝐇:\mathrm{\Omega }\mathrm{\Omega }`$ is an arbitrary symmetric (bounded) operator; the Hamiltonian equations have the form:
$$\dot{q}=𝐇_{21}q+𝐇_{22}p,\dot{p}=(𝐇_{11}q+𝐇_{12}p),$$
or
$$\dot{\psi }=\left(\begin{array}{cc}\dot{q}\hfill & \\ \dot{p}\hfill & \end{array}\right)=J𝐇\psi $$
(71)
Thus quadratic Hamilton functions induce linear Hamilton equations. From (71) we get
$$\psi (t)=U_t\psi ,\text{where}U_t=e^{J𝐇t}.$$
The map $`U_t\psi `$ is a linear Hamiltonian flow on the phase space $`\mathrm{\Omega }.`$ Let us consider a self-adjoint operator $`𝐇SO_H^{}(2)`$: $`𝐇=\left(\begin{array}{cc}R\hfill & T\hfill \\ T\hfill & R\hfill \end{array}\right).`$ This operator defines the quadratic Hamilton function
$$(q,p)=\frac{1}{2}[(Rp,p)+2(Tp,q)+(Rq,q)],$$
where the operator $`R`$ is symmetric and the operator $`T`$ is skew symmetric. Corresponding Hamiltonian equations have the form
$$\dot{q}=RpTq,\dot{p}=(Rq+Tp).$$
(72)
We point out that for a $`SO_H(2)`$-invariant Hamilton function, the Hamiltonian flow $`U_tSO_H^{}(2).`$ By considering the complex structure on the infinite-dimensional phase space $`\mathrm{\Omega }`$ we write the Hamiltonian equations (71) in the form of the Schödinger equation on $`H_c:`$
$$i\frac{d\psi }{dt}=𝐇\psi .$$
(73)
Its solution has the following complex representation: $`\psi (t)=U_t\psi ,U_t=e^{i𝐇t}.`$ We consider the Planck system of units in that $`h=1.`$ This is the complex representation of flows corresponding to quadratic $`SO_H(2)`$-invariant Hamilton functions.
By choosing $`H=L_2(𝐑^n)`$ we see that the interpretation of the solution of this equation coincides with the original interpretation of Schrödinger – this is a classical field
$$\psi (t,x)=(q(t,x),p(t,x).$$
Example 11.1. Let us consider an important class of Hamilton functions
$$(q,p)=\frac{1}{2}[(Rp,p)+(Rq,q)],$$
(74)
where $`R`$ is a symmetric operator. The corresponding Hamiltonian equations have the form:
$$\dot{q}=Rp,\dot{p}=Rq.$$
(75)
We now choose $`H=L_2(𝐑^3),`$ so $`q(x)`$ and $`p(x)`$ are components of the vector-field $`\psi (x)=(q(x),p(x)).`$ We can call fields $`q(x)`$ and $`p(x)`$ mutually inducing. The field $`p(x)`$ induces dynamics of the field $`q(x)`$ and vice versa, cf. with electric and magnetic components, $`q(x)=E(x)`$ and $`p(x)=B(x),`$ of the electromagnetic field, cf. Einstein and Infeld , p. 148: “Every change of an electric field produces a magnetic field; every change of this magnetic field produces an electric field; every change of …, and so on.” We can write the form (74) as
$$(q,p)=\frac{1}{2}_{𝐑^6}R(x,y)[q(x)q(y)+p(x)p(y)]𝑑x𝑑y$$
(76)
or
$$(\psi )=\frac{1}{2}_{𝐑^6}R(x,y)\psi (x)\overline{\psi }(y)𝑑x𝑑y,$$
(77)
where $`R(x,y)=R(y,x)`$ is in general a distribution on $`𝐑^6.`$ We call such a kernel $`R(x,y)`$ a self-interaction potential for the background field $`\psi (x)=(q(x),p(x)).`$We point out that $`R(x,y)`$ induces a self-interaction of each component of the $`\psi (x),`$ but there is no cross-interaction between components $`q(x)`$ and $`p(x)`$ of the vector-field $`\psi (x).`$
## 12 Invariant Gaussian measures of the Hamilton-Schrödinger dynamics and stationary pure states
All Gaussian measures considered in this section are supposed to be $`SO_H(2)`$-invariant. As we have seen, In our approach so called pure states $`\mathrm{\Psi },\mathrm{\Psi }=1,`$ are labels for Gaussian measures concentrated on one dimensional (complex) subspaces $`\mathrm{\Omega }_\mathrm{\Psi }`$ of the infinite-dimensional phase-space $`\mathrm{\Omega }.`$ In this section we study the case of so called stationary (pure) states in more detail. The $`\alpha `$-scaling does not play any role in present considerations. Therefore we shall not take it into account. We consider a pure state $`\mathrm{\Psi },\mathrm{\Psi }=1,`$ as the label for the Gaussian measure $`\nu _\mathrm{\Psi }`$ having the zero mean value and the complexification of the covariance operator
$$j(\rho _{\nu _\mathrm{\Psi }})=\mathrm{\Psi }\mathrm{\Psi }.$$
Theorem 12.1. Let $`\nu `$ be a Gaussian measure (with zero mean value) concentrated on the one-dimensional (complex) subspace corresponding to a normalized vector $`\mathrm{\Psi }`$. Then $`\nu `$ is invariant with respect to the unitary dynamics $`U_t=e^{it𝐇},`$ where $`𝐇:\mathrm{\Omega }\mathrm{\Omega }`$ is a bounded self-adjoint operator, iff $`\mathrm{\Psi }`$ is an eigenvector of $`𝐇.`$
Proof. A). Let $`𝐇\mathrm{\Psi }=\lambda \mathrm{\Psi }.`$ The Gaussian measure $`U_t^{}\nu `$ has the complexification of the covariance operator
$$j(\rho _t)=U_t(\mathrm{\Psi }\mathrm{\Psi })U_t^{}=U_t\mathrm{\Psi }U_t\mathrm{\Psi }=e^{it\lambda }\mathrm{\Psi }e^{it\lambda }\mathrm{\Psi }=\mathrm{\Psi }\mathrm{\Psi }.$$
Since all measures under consideration are Gaussian, this implies that $`U_t^{}\nu =\nu .`$ Thus $`\nu `$ is an invariant measure.
B). Let $`U_t^{}\nu =\nu `$ and $`\nu =\nu _\mathrm{\Psi }`$ for some $`\mathrm{\Psi },\mathrm{\Psi }=1.`$ We have that $`U_t\mathrm{\Psi }U_t\mathrm{\Psi }=\mathrm{\Psi }\mathrm{\Psi }.`$ Thus, for any $`\psi _1,\psi _2\mathrm{\Omega },`$ we have
$$\psi _1,U_t\mathrm{\Psi }U_t\mathrm{\Psi },\psi _2=\psi _1,\mathrm{\Psi }\mathrm{\Psi },\psi _2.$$
(78)
Let us set $`\psi _2=\mathrm{\Psi }.`$ We obtain: $`\psi _1,\overline{c(t)}U_t\mathrm{\Psi }=\psi _1,\mathrm{\Psi },`$ where $`c(t)=U_t\mathrm{\Psi },\mathrm{\Psi }.`$ Thus $`\overline{c(t)}U_t\mathrm{\Psi }=\mathrm{\Psi }.`$We point out that $`c(0)=\mathrm{\Psi }^2=1.`$ Thus $`\overline{c^{}(0)}\mathrm{\Psi }i𝐇\mathrm{\Psi }=0,`$ or $`𝐇\mathrm{\Psi }=i\overline{c^{}(0)}\mathrm{\Psi }.`$ Thus $`\mathrm{\Psi }`$ is an eigenvector of $`𝐇`$ with the eigenvalue $`i\overline{c^{}(0)}.`$ We remark that $`c^{}(0)=i𝐇\mathrm{\Psi },\mathrm{\Psi };`$ so $`\overline{c^{}(0)}=i𝐇\mathrm{\Psi },\mathrm{\Psi }.`$ Hence, $`\lambda =i\overline{c^{}(0)}=𝐇,\mathrm{\Psi },\mathrm{\Psi }.`$
Conclusion. In PCSFT stationary states of the quantum Hamiltonian (represented by a bounded self-adjoint operator $`𝐇)`$ are labels for Gaussian one-dimensional measures (with the zero mean value) that are invariant with respect to the Schrödinger dynamics $`U_t=e^{it𝐇}`$.
We now describe all possible Gaussian measures which are $`U_t`$-invariant.
Theorem 12.2. Let $`𝐇`$ be a bounded self-adjoint operator with purely discrete nondegenerate spectrum: $`𝐇\mathrm{\Psi }_k=\lambda _k\mathrm{\Psi }_k,`$ so $`\{\mathrm{\Psi }_k\}`$ is an orthonormal basis consisting of eigenvectors of $`𝐇.`$ Then any $`U_t`$-invariant Gaussian measure $`\nu `$ (with the zero mean value) has the complexification of the covariance operator:
$$j(\rho )=\underset{k=1}{\overset{\mathrm{}}{}}c_k\mathrm{\Psi }_k\mathrm{\Psi }_k,c_k0,$$
(79)
and vice versa.
Proof. A). Let $`j(\rho )`$ has the form (79). Then
$$j(\rho _{U_t^{}\nu })=U_tj(\rho )U_t^{}=\underset{k=1}{\overset{\mathrm{}}{}}c_ke^{i\lambda _kt}\mathrm{\Psi }_ke^{i\lambda _kt}\mathrm{\Psi }_k=j(\rho ).$$
(80)
Since measures are Gaussian, this implies that $`U_t^{}\nu =\nu `$ for any $`t.`$
B). Let $`U_t^{}\nu =\nu `$ for any $`t.`$ We remark that the complexification of any covariance operator $`\rho `$ can be represented in the form:
$$j(\rho )=\underset{k=1}{\overset{\mathrm{}}{}}j(\rho )\mathrm{\Psi }_k,\mathrm{\Psi }_k\mathrm{\Psi }_k\mathrm{\Psi }_k+\underset{kl}{}j(\rho )\mathrm{\Psi }_k,\mathrm{\Psi }_l\mathrm{\Psi }_k\mathrm{\Psi }_l.$$
(81)
We shall show that $`j(\rho )\mathrm{\Psi }_k,\mathrm{\Psi }_j=0`$ for $`kj.`$ Denote the operator corresponding to $`_{kj}`$ by Z. We have
$$U_tZU_t\psi _1,\psi _2=\underset{km}{}j(\rho )\mathrm{\Psi }_k,\mathrm{\Psi }_me^{it(\lambda _m\lambda _k)}\mathrm{\Psi }_k,\psi _2\psi _1,\mathrm{\Psi }_m=Z\psi _1,\psi _2.$$
(82)
Set $`\psi _1=\mathrm{\Psi }_j,\psi _2=\mathrm{\Psi }_k.`$ Then
$$U_tZU_t^{}\mathrm{\Psi }_m,\mathrm{\Psi }_k=j(\rho )\mathrm{\Psi }_k,\mathrm{\Psi }_me^{it(\lambda _m\lambda _k)}=j(\rho )\mathrm{\Psi }_k,\mathrm{\Psi }_m.$$
(83)
Thus $`j(\rho )\mathrm{\Psi }_k,\mathrm{\Psi }_m=0,km.`$
## 13 Stability of hydrogen atom
As we have seen, in PCSFT so called stationary (pure) states of quantum mechanics can play the role of labels for Gaussian measures (which are $`SO_H(2)`$-invariant and have zero mean value) that are $`U_t`$-invariant. We now apply our standard $`\alpha `$-scaling argument and we see that a stationary state $`\mathrm{\Psi }`$ is a label for the Gaussian measure $`\mu _\mathrm{\Psi }`$ with $`j(\rho _{\mu _\mathrm{\Psi }})=\alpha \mathrm{\Psi }\mathrm{\Psi }.`$ This measure is concentrated on one-dimensional (complex) subspace $`\mathrm{\Pi }_\mathrm{\Psi }`$ of phase space $`\mathrm{\Omega }.`$ Therefore each realization of an element of the Gaussian ensemble of classical fields corresponding to the statistical state $`\mu _\mathrm{\Psi }`$ gives us the field of the shape $`\mathrm{\Psi }(x),`$ but magnitudes of these fields vary from one realization to another. But by the well known Chebyshov inequality probability that $`(\mathrm{\Psi })=_{𝐑^3}|\mathrm{\Psi }(x)|^2𝑑x`$ is large is negligibly small.
Thus in the stationary state we have Gaussian fluctuations of very small magnitudes of the same shape $`\mathrm{\Psi }(x).`$ In PCSFT a stationary quantum state can not be identified with a stationary classical field, but only with an ensemble of fields having the same shape $`\mathrm{\Psi }(x).`$ Let us now compare descriptions of dynamics of electron in hydrogen atom given by quantum mechanics and our prequantum field theory.
In quantum mechanics stationary bound states of hydrogen atom are of the form:
$$\mathrm{\Psi }_{nlm}(r,\theta ,\varphi )=c_{n,l}R^lL_{n+l}^{2l+1}(R)e^{R/2}Y_l^m(\theta ,\varphi ),$$
(84)
where $`R=\frac{2r}{na_0},`$ and $`a_0=\frac{h^2}{\mu e^2}`$ is a characteristic length for the atom (Bohr radius). We are mainly interested in the presence of the component $`e^{R/2}.`$
In PCSFT this stationary bound state is nothing else, but the label for the Gaussian measure $`\rho _{\mathrm{\Psi }_{nlm}}`$ which is concentrated on the subspace $`\mathrm{\Omega }_{\mathrm{\Psi }_{nlm}}.`$ Thus PCSFT says that “electron in atom” is nothing else than Gaussian fluctuations of a certain classical field, namely the field $`\mathrm{\Psi }_{nlm}(r,\theta ,\varphi ):`$
$$\psi _{nlm}(r,\theta ,\varphi ;\psi )=\gamma (\psi )\mathrm{\Psi }_{nlm}(r,\theta ,\varphi ),$$
(85)
where $`\gamma (\psi )`$ is the C-valued Gaussian random variable: $`E\gamma =0,E|\gamma |^2=\alpha .`$
The intensity of the field $`\mathrm{\Psi }_{nlm}(r,\theta ,\varphi ,\psi )`$ varies, but the shape is the same. Therefore this random field does not produce any significant effect for large $`R`$ (since $`e^{R/2}`$ eliminates such effects).
Thus in PCSFT the hydrogen atom stable, since the prequantum random fields $`\psi _{nlm}(r,\theta ,\varphi ;\psi )`$ have a special shape (decreasing exponentially $`R\mathrm{}).`$
## 14 Appendixes
### 14.1 Classical representation for spin operators
The Pauli matrices are a set of $`2\times 2`$ complex Hermitian and unitary matrices. They are: $`\sigma _1=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right),`$ $`\sigma _2=\left(\begin{array}{cc}0& i\\ i& 0\end{array}\right),`$ $`\sigma _3=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right).`$ Let $`H_c=C^2`$ with the complex coordinates $`z=(z_1,z_2),z_j=q_j+ip_j,j=1,2,`$ and $`\mathrm{\Omega }=𝐑^2\times 𝐑^2`$ with the real coordinates $`\omega =(q_1,q_2,p_1,p_2).`$ We consider spin operators: $`\sigma (a)=_j^3a_j\sigma _j:𝐂^2𝐂^2,a=(a_1,a_2,a_3).`$ Let us consider real matrices $`\sigma _j^{(r)}=j^1(\sigma _j):`$
$$\sigma _1^{(r)}=\left(\begin{array}{cc}\sigma _1& 0\\ 0& \sigma _1\end{array}\right)=\left(\begin{array}{cccc}0& 1& 0& 0\\ 1& 0& 0& 0\\ 0& 0& 0& 1\\ 0& 0& 1& 0\end{array}\right),\sigma _2^{(r)}=\left(\begin{array}{cc}0& i\sigma _2\\ i\sigma _2& 0\end{array}\right)=\left(\begin{array}{cccc}0& 0& 0& 1\\ 0& 0& 1& 0\\ 0& 1& 0& 0\\ 1& 0& 0& 0\end{array}\right),$$
$$\sigma _3^{(r)}=\left(\begin{array}{cc}\sigma _3& 0\\ 0& \sigma _3\end{array}\right)=\left(\begin{array}{cccc}1& 0& 0& 0\\ 0& 1& 0& 0\\ 0& 0& 1& 0\\ 0& 0& 0& 1\end{array}\right).$$
We remark that these are not Dirac matrices. We set
$$\sigma ^{(r)}(a)j^1(\sigma (a))=\underset{j}{\overset{3}{}}a_j\sigma _j^{(r)}:𝐑^4𝐑^4,a=(a_1,a_2,a_3).$$
and consider classical random spin-variables: $`f_a(\omega )=\frac{1}{2}(\sigma ^{(r)}(a)\omega ,\omega ).`$ Then $`T(f_a)=\sigma (a)`$ and for any $`SO_2(2)`$-invariant Gaussian measure $`\mu `$ on $`\mathrm{\Omega }=𝐑^2\times 𝐑^2`$ with dispersion $`\alpha +O(\alpha )`$ we have: $`\frac{1}{\alpha }_{𝐑^4}f_a(\omega )𝑑\mu (\omega )=\mathrm{Tr}\mathrm{j}(\rho )\sigma (\mathrm{a})+\mathrm{O}(\alpha ),`$ where $`\rho `$ is the covariance operator of $`\sqrt{\alpha }`$-scaling of the Gaussian measure $`\mu .`$ For example,
$$\frac{1}{\alpha }_{𝐑^4}(q_1q_2+p_1p_2)𝑑\mu (q_1,q_2,p_1,p_2)=\mathrm{Tr}\mathrm{j}(\rho )\sigma _1+\mathrm{O}(\alpha ),$$
$$\frac{1}{\alpha }_{𝐑^4}(p_1q_2p_2q_1)𝑑\mu (q_1,q_2,p_1,p_2)=\mathrm{Tr}\mathrm{j}(\rho )\sigma _2+\mathrm{O}(\alpha ),$$
$$\frac{1}{\alpha }_{𝐑^4}(q_1^2q_2^2+p_1^2p_2^2)𝑑\mu (q_1,q_2,p_1,p_2)=\mathrm{Tr}\mathrm{j}(\rho )\sigma _3+\mathrm{O}(\alpha ).$$
We also have: $`\frac{1}{\alpha }_{𝐑^4}(q_1^2+q_2^2+p_1^2+p_2^2)𝑑\mu (q_1,q_2,p_1,p_2)=\mathrm{Tr}\mathrm{j}(\rho )\mathrm{I}+\mathrm{O}(\alpha )=\mathrm{Tr}\mathrm{j}(\rho )+\mathrm{O}(\alpha ).`$ By introducing vectors $`\omega _1=(q_1,p_1)`$ and $`\omega _2=(q_2,p_2)`$ we rewrite these asymptotic equalities in shorter way:
$$\frac{1}{\alpha }_{𝐑^4}(\omega _1,\omega _2)𝑑\mu (\omega _1,\omega _2)=\mathrm{Tr}\mathrm{j}(\rho )\sigma _1+\mathrm{O}(\alpha ),$$
$$\frac{1}{\alpha }_{𝐑^4}(J\omega _1,\omega _2)𝑑\mu (\omega _1,\omega _2)=\mathrm{Tr}\mathrm{j}(\rho )\sigma _2+\mathrm{O}(\alpha ),$$
where $`J`$ is the symplectic operator and, finally,
$$\frac{1}{\alpha }_{𝐑^4}(\omega _1^2\omega _2^2)𝑑\mu (\omega _1,\omega _2)=\mathrm{Tr}\mathrm{j}(\rho )\sigma _3+\mathrm{O}(\alpha ).$$
Let us now consider Gaussian measures on $`\mathrm{\Omega }=𝐑^4`$ corresponding to pure quantum states. These are singular Gaussian measures which are concentrated on $`SO_2(2)`$-invariant planes in $`𝐑^4.`$ To determine such a measure, we should find its covariation operator.
Proposition 14.1. Let $`\mathrm{\Psi }=u+iv,u=(u_1,u_2)𝐑^2,v=(v_1,v_2)𝐑^2`$ be a pure quantum state and let $`\rho _\mathrm{\Psi }=\mathrm{\Psi }\times \mathrm{\Psi }.`$ Then $`T^1(\rho _\mathrm{\Psi })=\mu _\mathrm{\Psi },`$ where the Gaussian measure $`\mu _\mathrm{\Psi }`$ has the covariation operator $`B_\mathrm{\Psi }=\alpha D_\mathrm{\Psi }`$ for
$$D_\mathrm{\Psi }=\left(\begin{array}{cccc}g_1^2& (g_1,g_2)& 0& (Jg_1,g_2)\\ \\ (g_1,g_2)& g_2^2& (g_1,Jg_2)& 0\\ \\ 0& (g_1,Jg_2)& g_1^2& (g_1,g_2)\\ \\ (Jg_1,g_2)& 0& (g_1,g_2)& g_2^2\end{array}\right).$$
Here $`g_1=(u_1,v_1)`$ and $`g_2=(u_2,v_2)`$ are variables which are conjugate to $`\omega _1=(q_1,p_1)`$ and $`\omega _2=(q_2,p_2).`$
Proof. The real space realization of $`\rho _\mathrm{\Psi }`$ is given by the operator:
$$j^1(\rho _\mathrm{\Psi })=\left(\begin{array}{cc}uu+vv\hfill & vuuv\hfill \\ uvvu\hfill & uu+vv\hfill \end{array}\right).$$
We have in the chosen system of coordinates on the phase space:
$$uu=\left(\begin{array}{cc}u_1^2& u_1u_2\\ u_1u_2& u_2^2\end{array}\right),vv=\left(\begin{array}{c}\\ v_1^2& v_1v_2\\ v_1v_2& v_2^2\end{array}\right),$$
$$vu=\left(\begin{array}{cc}u_1v_1& u_2v_1\\ u_1v_2& u_2v_2\end{array}\right),uv=\left(\begin{array}{cc}v_1u_1& v_2u_1\\ v_1u_2& v_2u_2\end{array}\right).$$
Hence:
$$uu+vv=\left(\begin{array}{cc}u_1^2+v_1^2& u_1u_2+v_1v_2\\ u_1u_2+v_1v_2& u_2^2+v_2^2\end{array}\right),$$
$$vuuv=\left(\begin{array}{cc}0& u_2v_1v_2u_1\\ u_1v_2v_1u_2& 0\end{array}\right).$$
To illustrate better correspondence between real and complex state spaces, we now show directly that $`j(D_\mathrm{\Psi })=\mathrm{\Psi }\mathrm{\Psi }`$ for $`D_\mathrm{\Psi }`$ given by this Proposition. We have
$$j(D_\mathrm{\Psi })=\left(\begin{array}{cc}g_1^2& (g_1,g_2)\\ (g_1,g_2)& g_2^2& \end{array}\right)+i\left(\begin{array}{cc}0& (Jg_1,g_2)\\ (g_1,Jg_2)& 0\end{array}\right).$$
This operator acts to a complex vector $`z=(z_1,z_2)`$ in the following way:
$$z_1^{}(j(D_\mathrm{\Psi })z)_1=g_1^2z_1+[(g_1,g_2)+i(Jg_1,g_2)]z_2,$$
$$z_2^{}(j(D_\mathrm{\Psi })z)_2=[(g_1,g_2)+i(g_1,Jg_2)]z_1+g_2^2z_2.$$
On the other hand, $`\mathrm{\Psi }\mathrm{\Psi }(z)=<z,\mathrm{\Psi }>\mathrm{\Psi }=(z_1\overline{\mathrm{\Psi }}_1+z_2\overline{\mathrm{\Psi }}_2)\mathrm{\Psi }.`$ Here
$$z_1^{}=(u_1iv_1)(u_1+iv_1)z_1+(u_2iv_2)(u_1+iv_1)z_2,$$
$$z_2^{}=(u_1iv_1)(u_2+iv_2)z_1+(u_2iv_2)(u_2+iv_2)z_2.$$
Thus
$$z_1^{}=(u_1^2+v_1^2)z_1+[(u_1u_2+v_1v_2)+i(u_2v_1u_1v_2)]z_2,$$
$$z_2^{}=(u_2^2+v_2^2)z_2+[(u_1u_2+v_1v_2)+i(u_1v_2u_2v_1)]z_2.$$
Let us consider Gaussian measures corresponding to pure states for spin up and spin down, $`|1>=\left(\begin{array}{c}1\\ 0\end{array}\right)`$ and $`|0>=\left(\begin{array}{c}0\\ 1\end{array}\right).`$ For vector $`|1>,`$ we have: $`u_1=1,u_2=v_1=v_2=0.`$ Thus $`D_{|1>}=\left(\begin{array}{cccc}1& 0& 0& 0\\ 0& 0& 0& 0\\ 0& 0& 1& 0\\ 0& 0& 0& 0\end{array}\right)`$ and the Fourier transform of the measure $`\mu _{|1>}`$ is given by:
$$\stackrel{~}{\mu }_{|1>}(\xi _1,\xi _2,\eta _1,\eta _2)=e^{\frac{\alpha }{2}(\xi _1^2+\eta _1^2)}.$$
This is the standard Gaussian measure on the plane $`q_2=0,p_2=0`$ having the density: $`d\mu _{|1>}(q_1,p_1)=\frac{1}{2\pi \alpha }e^{\frac{1}{2\alpha }(q_1^2+p_1^2)}.`$ In the same way $`\mu _{|0>}`$ is the standard Gaussian measure on the plane $`q_1=0,p_1=0`$ having the density: $`d\mu _{|0>}(q_2,p_2)=\frac{1}{2\pi \alpha }e^{\frac{1}{2\alpha }(q_2^2+p_2^2)}.`$ Let us now consider the Gaussian measure corresponding to superposition of spin up and spin down states:
$$\mathrm{\Psi }_\theta =\frac{1}{\sqrt{2}}(|0>+e^{i\theta }|1>).$$
Here $`u_1=\mathrm{cos}\theta ,v_1=\mathrm{sin}\theta .`$ Hence
$$D_{\mathrm{\Psi }_\theta }=\left(\begin{array}{cccc}\mathrm{cos}^2\theta & \mathrm{cos}\theta \mathrm{sin}\theta & 0& 0\\ \mathrm{cos}\theta \mathrm{sin}\theta & \mathrm{sin}^2\theta & 0& 0\\ 0& 0& \mathrm{cos}^2\theta & \mathrm{cos}\theta \mathrm{sin}\theta \\ 0& 0& \mathrm{cos}\theta \mathrm{sin}\theta & \mathrm{sin}^2\theta \end{array}\right)$$
and the Fourier transform of $`\mu _{\mathrm{\Psi }_\theta }`$ is given by
$$\stackrel{~}{\mu }_{\mathrm{\Psi }_\theta }(\xi _1,\xi _2,\eta _1,\eta _2)=e^{\frac{\alpha }{2}[(\mathrm{cos}\theta \xi _1+\mathrm{sin}\theta \xi _2)^2+(\mathrm{cos}\theta \eta _1+\mathrm{sin}\theta \eta _2)^2]}.$$
Thus pure states $`\mathrm{\Psi }_\theta `$ correspond to the standard Gaussian measures concentrated on planes obtained by rotations.
### 14.2 Comparing with no-go theorems of von Neumann, Cohen-Specker and Bell
There are no-go theorems for mathematical attempts to have a map from classical variables to quantum operators which preserves statistics, e.g., theorems of von Neumann, Cohen-Specker and Bell, see , . The no-go theorems say: No such map exists. In this paper we constructed such a map. What goes?
Our construction does not contradict to known no-go theorems, since our map $`T`$ does not satisfy some conditions of those theorems. An important condition in all such theorems is that the range of values of a classical variable $`f`$ should coincide with the spectrum of the corresponding quantum operator $`T(f)`$ – “the range of values postulate.” This postulate is violated in our framework. As we have seen, the classical spin variables are continuous and the quantum spin operators have discrete spectrum. Nevertheless, classical averages can be approximated by quantum. Our prequantum classical statistical model is not about observations, but about ontic reality (reality as it is when nobody looks at it).
Henry Stapp pointed out : “The problem, basically, is that to apply quantum theory, one must divide the fundamentally undefined physical world into two idealized parts, the observed and observing system, but the theory gives no adequate description of connection between these two parts. The probability function is a function of degrees of freedom of the microscopic observed system, whereas the probabilities it defines are probabilities of responses of macroscopic measuring devices, and these responses are described in terms of quite different degrees of freedom.” Since we do know yet from physics so much about features of classical $``$ quantum correspondence map $`T,`$ we have the freedom to change some conditions which were postulated in the known no-go theorems – for example, the range of values condition. Rejection of this assumption is quite natural, since, as was pointed by Stapp, a classical variable $`f`$ and its quantum counterpart $`T(f)`$ depend on completely different degrees of freedom.
### 14.3 Is prequantum classical statistical field theory nonlocal?
As we have seen, PCSFT does not contradict to the known no-go theorems, in particular, to Bell’s theorem. Therefore this theory might be local. However, it is not easy to formulate the problem of locality/nonlocality in the PCSTF-framework. It is not about observations. Thus we could not apply Bell’s approach to locality as locality of observations. On the other hand, on the ontic level PCSTF operates not with particles, but with fields. At the first sight, such a theory is nonlocal by its definition, since fields are not localized. But in field theory there was established a different viewpoint to locality and we know that both classical and quantum field theories are local. To formulate the problem of locality for PCSTF, we should proceed in the same way. Therefore we should develop a relativistic version of PCSTF. There are some technical and even ideological problems. As we know, relativistic quantum mechanics is not a well established theory (at least this is a rather common opinion). Thus it is meaningless to develop a relativistic variant of PCSTF which would reproduce relativistic quantum mechanics. The most natural way of development is to construct a kind of PCSTF not for quantum mechanics, but for quantum field theory and study the problem of locality in such a framework. It is an interesting and complicated problem which will be studied in coming papers of the author. |
warning/0506/astro-ph0506077.html | ar5iv | text | # Infrared Surface Brightness Distances to Cepheids: a comparison of Bayesian and linear-bisector calculations
## 1 Introduction
The infrared surface brightness technique is a powerful method for determining distances to Cepheid variables (Welch 1994, Fouqué & Gieren 1997). It is independent of other astrophysical distance scales, nearly independent of errors in reddening, and may be applied to arbitrarily-chosen, individual Cepheids at great distances. However, the implementation of the infrared surface brightness technique and its predecessor, the visual surface brightness method, have been criticized as not mathematically rigorous in their solutions to the surface brightness equations (Laney & Stobie 1995, Barnes & Jefferys 1999). This leads to the possibility that the distances, radii, and their uncertainties may be erroneous.
This is of current interest as Storm et al. (2004) have found a short distance to the LMC based on a linear-bisector analysis of six Cepheids in the LMC cluster NGC1866. They also found a slope to the LMC Cepheid period-luminosity relation that is substantially different than found in the OGLE magnitudes (Udalski et al. 1999). It is important to determine if their results are affected by the mathematical method used in the surface brightness analysis.
To address the larger issue Barnes et al. (2003) developed a Bayesian Markov Chain Monte Carlo (MCMC) solution to the surface brightness equations that is mathematically rigorous. In the current paper we do a direct comparison for a significant sample of Cepheids between the Bayesian MCMC solution and recent linear-bisector calculations (Storm et al. 2004, Gieren et al. 2005) to explore possible differences. The two methods used identical data, identical surface brightness equations and identical physical constants to ensure that only the mathematical approaches were compared.
In the next section we introduce the surface brightness method for determining distances and radii and present the background of the infrared surface brightness method. We then discuss the data that are used for the two calculations. We review the linear-bisector calculations and the Bayesian MCMC calculations. In section 6 we compare the results of the two methods for 38 Galactic Cepheids. Finally we discuss the importance of the agreement and the differences that we find.
## 2 The Surface Brightness Method
### 2.1 The surface brightness equations
The infrared surface brightness method is a modification of the visual surface brightness technique developed by Barnes et al. (1977) and thus shares the same computational algorithm. Because solution of the equations is the important issue, we introduce the equations in some detail. Useful discussions of previous work have been given by Gieren, Barnes & Moffett (1993), Fouqué & Gieren (1997), Nordgren et al. (2002), Fouqué, Storm, & Gieren (2003), and Barnes et al. (2003).
Barnes & Evans (1976) and Barnes, Evans, & Parsons (1976) defined a visual surface brightness parameter $`F_V`$ as
$$F_V=4.22070.1V{}_{0}{}^{}0.5\mathrm{log}\varphi $$
(1)
and also as,
$$F_V=\mathrm{log}T_e+0.1BC,$$
(2)
where $`V_0`$ is the stellar visual magnitude corrected for interstellar extinction, $`\varphi `$ is the stellar angular diameter expressed in milliarcseconds, $`T_e`$ is the effective temperature and $`BC`$ is the bolometric correction.
They demonstrated that $`F_V`$ is well correlated with Johnson color index $`(VR)_0`$ for a very wide range of stellar types.
$$F_V=A+B(VR)_0$$
(3)
equation (3) is called the visual surface brightness relation.
These relations led Barnes et al. (1977) to infer a distance scale for Cepheids as follows. At each time $`t`$ in the pulsation of the Cepheid, equations (1) and (3) may be combined to obtain the angular diameter variation of the star, $`\varphi (t)`$,
$$4.22070.1V{}_{0}{}^{}(t)0.5\mathrm{log}\varphi (t)=A+B(VR){}_{0}{}^{}(t)$$
(4)
In addition we infer the Cepheid’s linear radius variation $`\mathrm{\Delta }R(t)`$ about the mean radius from an integration of the radial velocity curve, $`V_r(t)`$,
$$\mathrm{\Delta }R(t)=p\left(V_r(t)V_\gamma \right)𝑑t$$
(5)
where the factor $`p`$ converts observed radial velocity to the star’s pulsational velocity and $`V_\gamma `$ is the center-of-mass radial velocity of the star. Integration of the discrete radial velocity data requires that the velocity variation be appropriately modeled.
Substituting into equation (4), the relation among mean angular diameter $`\varphi _0`$, linear diameter $`\mathrm{\Delta }R(t)`$, and distance $`r`$, we obtain
$$4.22070.1V{}_{0}{}^{}(t)0.5\mathrm{log}(\varphi _0+2000\mathrm{\Delta }R(t)/r)=A+B(VR){}_{0}{}^{}(t).$$
(6)
where $`\varphi _0`$ is in milliarcseconds, $`r`$ is in parsecs, and $`\mathrm{\Delta }R(t)`$ is in AU. The factor $`2000`$ converts radius to diameter and arcseconds to milliarcseconds.
Until very recently direct solution to equation (6) for distance and diameter has proved daunting. Various methods have been tried, but for reasons given by Barnes & Jefferys (1999), none of the methods is rigorous.
Barnes et al. (1977), Gieren et al. (1993, 1997) and Welch (1994) simplified the problem by solving equation (4) for $`\varphi (t)`$, solving equation (5) for $`\mathrm{\Delta }R(t)`$ and then solving for $`r`$ and $`\varphi _0`$ using ordinary least-squares solution on equation (7) taking $`\mathrm{\Delta }R(t)`$ as the independent variable.
$$\mathrm{\Delta }R(t)=r(\varphi _0+\varphi (t))/2000$$
(7)
However, the least-squares calculations do not properly treat the errors-in-variables problem that arises from uncertainty in both $`\mathrm{\Delta }R(t)`$ and $`\varphi (t)`$. The best that can be done using least-squares is the linear-bisector solution. Isobe et al. (1990) showed that linear-bisector performs better than other least-squares solutions when the problem is symmetric in the variables, as is the case here. Storm et al. (2004) adopted the linear-bisector method.
Following Balona (1977), Laney & Stobie (1995) linearized equation (6) and applied a maximum-likelihood calculation to its solution. They used an iterative maximum-likelihood method to solve for the larger amplitude Cepheids for which the linearization is invalid. Their calculation addresses the errors-in-variables problem. However, results from the maximum-likelihood method can be quite sensitive to accurate knowledge of the uncertainties in the data, as discussed by Gieren et al. 1997.
All the above methods require that the observed radial velocities be modeled in order to do the integration of equation (5). This creates a model selection problem. Most researchers choose a Fourier series to fit the radial velocities, but the number of terms to include in the series is subjectively chosen.
Finally, whether least-squares or maximum-likelihood, the above calculations fail to treat properly the propagation of observational error through the radial velocity integral, equation (5). Balona (1977) introduced a widely used approximation to the error in $`\mathrm{\Delta }R(t)`$ based on the uncertainty in the radial velocity data and on the assumption that the radial velocity data are equally spaced in pulsation phase, which is rarely true.
It was to address these issues that the Bayesian MCMC method was developed by Barnes et al. (2003). Their calculation correctly solves the errors-in-variables problem, objectively selects the number of terms in the Fourier series fits, and correctly propagates the observational error through the radial velocity integration. Moreover, the Bayesian MCMC calculation does not demand that a particular Fourier series fit the data; rather, each Fourier series that is fit to the data has a particular posterior probability. That probability is used as a weight in determining the other quantities sought in the solution, i.e., distance, mean radius, etc.
A large number of calculations of variable star distances and radii by surface brightness methods are present in the literature. It is therefore of interest to determine the extent to which the above listed deficiencies affect those results. We will do this by comparing the linear-bisector calculations with Bayesian MCMC calculations.
### 2.2 The infrared surface brightness technique
Welch (1994) first showed that use of the infrared combination $`K{}_{0}{}^{},(VK)_0`$ in place of $`V{}_{0}{}^{},(VR)_0`$ in equation (6) has significant advantage in the precision of the distances and radii for Cepheids. He found an improvement by a factor of three in the distance uncertainty for the Cepheid U Sgr. Welch attributed the improved precision to several factors. First, the color index $`(VK)_0`$ is as good an indicator of surface brightness as bluer color indices but much less sensitive to the complications of line blanketing and surface gravity. Second, the $`K_0`$ magnitude lies on the Rayleigh-Jeans tail of the flux distribution and is therefore less affected by the surface brightness variation. It is obvious that a magnitude more sensitive to the radius variation that is then corrected for surface brightness variation using a color index that is a more accurate indicator of surface brightness ought to yield superior results.
Laney & Stobie (1995) independently examined two optical magnitude-color index combinations and two infrared combinations to determine which would give the most precise radii of Cepheids. From both model atmosphere considerations and the precision of 49 Cepheid calculations, they concluded that $`K{}_{0}{}^{},(VK)_0`$ and $`K{}_{0}{}^{},(VJ)_0`$ were superior to the optical indices. (Examination of their Tables 5–6 suggests that $`K{}_{0}{}^{},(VK)_0`$ is slightly better, confirming Welch’s choice.) Laney & Stobie also noted that the presence of a companion to the Cepheid would have less effect on infrared indices than on optical ones, as Cepheids are more likely to have a blue main-sequence companion than a red giant companion from stellar evolution considerations.
Extending Welch’s (1994) work on U Sgr, Fouqué & Gieren (1997) compared distance and radius results using $`V{}_{0}{}^{},(VR)_0`$; $`V{}_{0}{}^{},(VK)_0`$; and $`K{}_{0}{}^{},(JK)_0`$ with a somewhat improved data base and improved surface brightness equations. They found the same distances and radii from all three combinations, but much superior precision for the infrared color indices. The combination $`V{}_{0}{}^{},(VK)_0`$ seemed to be slightly preferred on the basis of precision. When adjusted to the same surface brightness equations, Fouqué & Gieren’s results agree within the errors with those of Welch and of Laney & Stobie.
With a sample of 16 Galactic cluster Cepheids, Gieren, Fouqué, & Gómez (1997) repeated the comparison of the above three combinations. Again they concluded that the infrared color indices are superior to the optical in precision and also agree well with each other. On the other hand, the $`V{}_{0}{}^{},(VR)_0`$ solutions gave distances and radii $`14\%`$ larger than the infrared solutions. While they indicated no preference between $`V{}_{0}{}^{},(VK)_0`$ and $`K{}_{0}{}^{},(JK)_0`$, their Table 3 shows that $`V{}_{0}{}^{},(VK)_0`$ gave better precision than $`K{}_{0}{}^{},(JK)_0`$ in 13 of 16 cases.
Based on these studies, we adopt the combination $`V{}_{0}{}^{},(VK)_0`$ as the basis for our comparison of the mathematical methods.
## 3 The Data
Storm et al. (2004) analyzed $`34`$ Galactic Cepheids for distances and radii using the $`V{}_{0}{}^{},(VK)_0`$ infrared surface brightness technique. Gieren et al. (2005) enlarged the sample by adding four more stars. We have used the full set of 38 Cepheids. The infrared surface brightness relation adopted in both those studies and by us is the one determined by Fouqué & Gieren (1997):
$$F_V=3.9470.131(VK)_0$$
(8)
The individual stellar data required for the analyses are the photometric measures $`V,(VK)`$, reddening $`E(BV)`$, radial velocities $`V_r`$, pulsation period $`P`$, and pulsation phases $`\theta `$. Storm et al. list the sources for photometry and radial velocities in their Table 1. To ensure that we used identical data in both calculations, the Bayesian analysis used the same input data files as used in the linear-bisector analysis.
We added Z Lac, Y Oph, S Sge, and CS Vel to the program using data referenced in Gieren et al. (2005). In addition, the results for $`\mathrm{}`$ Car given by Storm et al. were revised to incorporate new radial velocities, which are also referenced in Gieren et al. (2005). Again, the Bayesian analysis used the same input data files as used in the linear-bisector analysis by Gieren et al.
Although the infrared surface brightness method is largely independent of the interstellar extinction, it is important that the same extinction and reddening be used in our two calculations. We follow the earlier studies and adopt $`E(BV)`$ from Fernie’s (1990) tabulation for Cepheids. These $`E(BV)`$ values are listed in Table 1 along with the periods of the Cepheids. We adopted $`A_V=3.26E(BV)`$ and $`E(VK)=2.88E(BV)`$, slightly different than the reddening law used by Storm et al. The linear-bisector calculations were repeated with the new reddening law for all 38 Cepheids.
In the integration of the radial velocity curve to obtain the linear displacements, equation (5), the value of $`p`$ is required. Some studies have adopted a single value of $`p`$ for all Cepheids, others have adopted individual values. This choice has no effect upon our comparison of linear-bisector and Bayesian computations, provided the same value is used in both calculations for a specific Cepheid. Storm et al. and Gieren et al. used a relation between $`p`$ and period developed by Gieren et al. (1989) to approximate model atmosphere results by Hindsley & Bell (1986):
$$p=1.390.03logP$$
(9)
where $`P`$ is the pulsation period. We adopt this relation as well, using the periods given in Table 1 for the Bayesian MCMC calculation.
## 4 The Linear-bisector Computations
Storm et al. (2004) and Gieren et al. (2005) applied a linear-bisector, least-squares solution to the surface brightness equations. They solved equation (4) for the angular diameter variation $`\varphi (t)`$. The radial velocity data $`V_r(t)`$ were linearly interpolated in phase between the observed velocities and equation (5) integrated in 0.01 steps in phase to obtain the $`\mathrm{\Delta }R(t)`$ variation. Finally, equation (7) was solved for $`r`$ and $`\varphi _0`$.
To verify the solution for each Cepheid, the angular diameter variation and the displacement variation were plotted against phase in the pulsation cycle. For many of the stars, the angular diameter variation was a very poor match to the displacement curve in the phase interval $`0.81.0`$, as illustrated in Figure 1. Storm et al. discussed the source of the poor fit without coming to a conclusion and decided that the best course of action was to ignore this phase interval in the fit. They deleted the phase interval $`0.81.0`$ for all stars in their fit of equation (7) for $`r`$ and $`\varphi _0`$.
Some Cepheids showed a small phase shift between the photometry and the radial velocities. The phase shift causes a loop in the upper panel and a displacement between the fitted curve and the points in the lower panel. For the stars showing this effect a phase shift was determined by minimizing the scatter in the upper panel. The phase shift has already been imposed between the photometry and radial velocities used for Fig. 1 and thus these effects are not seen. The adopted phase shifts are listed in Table 1.
There are two constants adopted in a surface brightness calculation that must be the same in the linear-bisector and Bayesian calculations. We chose to fix the constants in the Bayesian calculation to those used in the Storm et al. paper. These constants are the constant term in the surface brightness definition, equation (1), and the conversion factor from angular diameter (in milliarcseconds) and distance (parsecs) to linear radius (solar radii). The first of these was taken to be 4.2207 and the second as 0.10727 solar radii per mas-parsec (see eq. (7)).
To perform the calculations Storm et al. used the FORTRAN subroutine SIXLIN which is available from Isobe et al. (1990). The computations were run on a Linux personal computer and took a fraction of a second per star.
The quantities determined in the linear-bisector calculation that are of interest to us here are the distance $`r`$, the mean linear radius $`R`$, computed from the distance and mean angular diameter $`\varphi _0`$, and their 1 $`\sigma `$ uncertainties. These quantities are given in Table 1.
## 5 The Bayesian Markov Chain Monte Carlo Computations
The Bayesian MCMC method that we applied here is described in detail by Barnes et al. (2003). It is important to realize that the Bayesian calculation treats the unknowns in the problem as probability distributions, not as specific values to be determined. The goal of the analysis is to determine the probability distribution for each parameter of interest, from which inferences may be drawn by appropriate means and variances.
To compute the posterior probability distribution for each parameter requires the likelihood function on a specific model, appropriate priors, and sampling strategies for all parameters. The likelihood function is the probability of obtaining the particular data given the model. Bayesian statistics encapsulate our understanding of the parameters in the model, prior to considering the data, in prior distributions. The posterior probability distribution is the product of the prior and the likelihood, appropriately normalized.
The full posterior probability distribution in this problem requires solution to integrals in the normalization that cannot be done analytically. However, the unnormalized posterior probability distribution avoids these integrals and can be used to generate a Monte Carlo sample from the full posterior probability distribution. The techniques used to generate the sample are Markov Chain techniques.
The model for the infrared surface brightness calculation is developed by first substituting the infrared color index for the visual color index in equation (6) and then rearrange as follows
$$(VK){}_{0}{}^{}(t)=\frac{1}{B}(4.22070.1V{}_{0}{}^{}(t)A0.5\mathrm{log}(\varphi _0+2000\mathrm{\Delta }R(t)\pi ))$$
(10)
where $`A,B`$ take the values given in equation (8) and where we have replaced $`1/r`$ with $`\pi `$, the parallax in arcseconds. Within this model the likelihood function is specified in a straightforward way and is given in equation (12) of Barnes et al. (2003). We model the photometry and the radial velocity data as drawn from normal distributions with variances given by the observational uncertainties. Because we do not trust the quoted observational uncertainties, we introduce a hyper-parameter scale factor on each variance to model deviation in the scatter from that expected from the quoted uncertainties. The time variations of the photometry and radial velocities are modeled by Fourier series of unknown order on the pulsation phase.
Barnes et al. discuss the priors adopted for each parameter of interest, but only one is relevant here. It is well-known that Cepheids are distributed within the plane of the Galaxy with an exponential decrease in density away from the plane. We adopted a prior on distance that reflects the flattened density distribution with a scale height of $`70\pm 10`$ pc. The results are insensitive to reasonable changes in the scale height.
The sampling strategy employed in this work is Markov Chain Monte Carlo using the Metropolis-Hastings and Gibbs algorithms, as described in Barnes et al. The art in this approach is to find sampling methods that explore the posterior probability distributions fully and efficiently. Internal tests provide guidance on the completeness and efficiency of the sampling. All the results presented here passed those tests. In the customary manner, we chose a burn-in phase to improve the model selection efficiency.
To be consistent with the linear-bisector calculation, we deleted from the Bayesian solution the photometry and radial velocities in pulsation phase interval $`0.81.0`$. In the general model, we allow for an unknown phase shift between the photometry and the radial velocities; in this calculation, we fixed the phase shift at the value determined by the linear-bisector calculation.
The calculations were run on a 1 GHz Macintosh G4 computer under system MacOS 10.3 using the statistical language R-1.8.0$`\beta `$ distributed by the R Development Core Team.<sup>1</sup><sup>1</sup>1See http://www.r-project.org/ . (The R language is available for other platforms.) We used a burn-in of 1,000 samples followed by 10,000 samples. Because the R code is interpreted code, the calculations run slowly; the computations here typically took about an hour per star.
Posterior probability distributions were determined for all the model parameters of interest. These are the parallax $`\pi `$, the mean angular diameter $`\varphi _0`$, the orders of the Fourier series on the apparent magnitudes and the radial velocities, the aforementioned hyper-parameters on the observational uncertainties, the mean $`V_0`$ magnitude (both intensity mean and magnitude mean), and the center-of-mass radial velocity. For our purpose here, only the parallax and mean angular diameter are important. These were converted to distance and mean linear radius within the code and listed with their uncertainties in Table 1. (The erratum published by Barnes et al. (2003) was addressed in the computation of the radii.)
## 6 Comparison of the Results
### 6.1 The distances and radii
Our goal is to determine differences in the results of the two calculations with respect to distance and radius. Figures 2 and 3 show the Bayesian values plotted against the linear-bisector values. There is no obvious difference between the values from the two calculations. To look for subtle differences, we computed the ratio of the Bayesian distance to the linear-bisector distance. A weighted least-squares fit of this ratio against $`log(P)`$ gives (Figure 4)
$$Ratio(r)=1.016(\pm 0.007)0.012(\pm 0.022)(logP1.113)$$
(11)
where the fit is centered on the mean period. A weighted fit for the radius ratio gives (Figure 5)
$$Ratio(R)=1.012(\pm 0.007)0.004(\pm 0.023)(logP1.113)$$
(12)
These show no evidence for any dependence on pulsation period. Therefore we take a weighted mean for each ratio: the distance calculations differ by $`1.5\%\pm 0.6\%`$, with the Bayesian being larger, and the radius calculations differ by $`1.1\%\pm 0.7\%`$, with the Bayesian again being larger. For comparison, the best individual distance and radius measurements in our dataset are for X Cyg, for which the Bayesian uncertainties are both $`\pm 2.6\%`$, larger than the possible systematic difference between the Bayesian and bisector calculations. Our first result is that the distances and radii computed by the two methods agree quite well.
### 6.2 Uncertainties in distance and radius
In Figures 4 and 5 the uncertainties can be seen to increase as the period is shorter. This is likely a result of the smaller pulsation amplitudes at shorter periods, which result in the photometric and velocity uncertainties having greater effect on the computed distances and radii.
More importantly, a glance at Table 1, or Figures 4 and 5, shows that the uncertainties in the distances and radii disagree substantially between the two calculations. Typically the Bayesian uncertainty is more than three times the linear-bisector uncertainty.
Because the concept of an uncertainty estimated from the Bayesian MCMC posterior probability distribution may not be clear, we show an example in Figure 6, the posterior probability distribution for the distance to U Sgr, which star is typical of our results. Over-plotted on the probability distribution is a normal distribution constructed for the same distance (592 pc), sigma ($`\pm 21`$ pc) and area (10,000 samples). The normal distribution describes the posterior probability distribution for the distance very well. This justifies our adopting the sigma of the corresponding normal distribution as a 1$`\sigma `$ estimator for the uncertainty in the distance (and similarly for the radius) determined in the Bayesian calculation. This estimator is then compared to the 1$`\sigma `$ estimator from the linear-bisector computation.
We add a caveat to the previous paragraph. Because our computation determines the stellar parallax, not the stellar distance, (see equation (10)) the posterior probability distribution for the distance can become asymmetric when the errors are large. As the uncertainties become large, the parallax posterior probability distribution becomes broad (large sigma). Its reciprocal, the distance posterior probability distribution, will also become broad and necessarily asymmetric to larger distances. The same asymmetry will arise for the radius posterior probability distribution in such cases because that distribution is the product of the angular diameter posterior probability distribution (symmetric) and the distance posterior probability distribution (asymmetric). Note that these asymmetric distributions are real and not mathematical artifacts; they properly represent our knowledge of the distance and radius, which is not true for least-squares or maximum-likelihood calculations on the same data. The latter methods assume symmetric errors by their very nature. Because this situation prevails for only a few stars in this sample, and only for stars with large errors, it has little effect on the weighted mean ratios of distances and radii quoted in the previous section.
As we did with the distances and radii themselves, we begin by examining the behavior of the ratio of the Bayesian uncertainty to the linear-bisector uncertainty for the same Cepheid. In Figures 7 and 8 we show these ratios for the distance and radius uncertainties plotted against $`logP`$. Unweighted least-squares fits in these figures yield
$$Ratio(r)=3.396(\pm 0.183)0.748(\pm 0.549)(logP1.113)$$
(13)
$$Ratio(R)=3.281(\pm 0.202)0.673(\pm 0.606)(logP1.113)$$
(14)
There is no apparent dependence of these ratios on pulsation period. Plots of the ratios of the distance uncertainties against distance and of the ratios of the radius uncertainties against radius are similarly uninformative.
The two ratios are, however, highly correlated with each other ($`R=0.99`$) as shown in Figure 9. Thus the underlying cause of the larger uncertainties in the Bayesian calculation is likely to be the same for the distance uncertainty and radius uncertainty.
In section 2.1 we noted that the linear-bisector calculation does not treat the errors-in-variables problem rigorously nor does it properly propagate uncertainty through the radial velocity integration. The second of these issues will certainly lead to an underestimate of the uncertainties in the computed distances and radii. Because the Bayesian MCMC calculation does correctly address these two computational issues, we interpret the large ratio of Bayesian to bisector uncertainty as measuring the amount by which the linear-bisector errors have been underestimated. This interpretation is supported by the fact that none of the linear-bisector uncertainties is larger than its Bayesian counterpart. Our second result is that the linear-bisector calculation underestimates the uncertainties in distance and in radius substantially, amounting to factors of 1.4–6.7 for this dataset. This large range implies that the ratio that is obtained depends on the specifics of the data for the Cepheid which varies from star to star.
## 7 Discussion
We set out to determine whether infrared surface brightness estimates of Cepheid distances and radii by the linear-bisector calculation are affected by the known mathematical shortcomings of that calculation. Based on comparison of Bayesian MCMC and linear-bisector calculations for 38 Cepheids using the same data, same surface brightness equations, and same physical constants, we find that the distances and radii are not adversely affected but that the uncertainties in these quantities are seriously underestimated in the linear-bisector calculation.
We find that Cepheid distances determined by the two calculations agree to $`1.5\%\pm 0.6\%`$ with the Bayesian distances being larger. This may be compared to the smallest individual uncertainty in distance found in the Bayesian calculation of $`\pm 2.6\%`$. Similarly we find that Cepheid radii determined by the two calculations agree to $`1.1\%\pm 0.7\%`$ with the Bayesian radii being larger. This may be compared to the smallest individual uncertainty in radius found in the Bayesian calculation of $`\pm 2.6\%`$. Any systematic difference between the mathematical approaches is both smaller than a $`2\sigma `$ effect and smaller than the typical single-star (Bayesian) uncertainty dictated by the data.
These results have an impact on interpretation of infrared surface brightness results for Cepheids. For example, Storm et al. (2005) used six Cepheids in an LMC cluster to infer a distance to the LMC of $`(mM)_0=18.30\pm 0.07`$ mag. by means of a linear-bisector solution for the infrared surface brightness equations. This is less than the generally accepted distance of 18.50 mag. From the present work, we can say that the linear-bisector calculation used by Storm et al. is not the cause of the smaller distance modulus.
Similarly, that work found a much smaller slope for the Cepheid PL relation in the LMC than had been found by the OGLE project. From the absence of a period dependence between the Bayesian MCMC distances and the linear bisector distances (Fig. 4) and from the overall agreement between the two, we can be certain that the smaller slope found by Storm et al. (2005) is not a result of using the linear-bisector, least-squares method.
The uncertainties in distance and radius determined by the Bayesian MCMC calculation are much larger than determined by the linear-bisector calculation. Given the known problems in a least-squares solution to the surface brightness equations, we interpret this as measuring the amount by which the linear-bisector computation underestimates the uncertainties. It is important to note that the ratio of Bayesian to linear-bisector uncertainty ranges from 1.4 to 6.7 in this set of 38 Cepheids. Clearly the amount by which the linear-bisector method underestimates the uncertainty depends on the specific nature of the data. This is expected, but inconvenient, as it is not possible to simply multiply published linear-bisector uncertainties by a constant correction factor.
Our result on the underestimation of the uncertainties in a linear-bisector calculation is supported in the previously mentioned paper by Storm et al. (2005). As the six Cepheids studied by them are in an LMC cluster, we can be confident that they are at the same distance. The scatter in their distances is then an estimate of the true uncertainty in the linear-bisector distance to the cluster. This scatter was found by Storm et al. (2005) to be twice the formal errors of the linear-bisector distances to the Cepheids, within the range of results determined here.
Other mathematical approaches have been used to solve the surface brightness equations for distance and radius. Ordinary least-squares assumes no error on one variable and all errors on the other. Inverse fit least-squares assumes the reverse. The limitations of ordinary linear least-squares solutions (direct and inverse) include not only underestimation of the errors, but also possible systematic bias in the resulting distances and radii as discussed by Laney & Stobie (1995) and Gieren et al. (1997).
The linear-bisector, least-squares calculation achieves a solution in-between the two other least-squares calculations (direct and inverse). Moreover, the linear-bisector error bar roughly corresponds to the difference between the two least-squares solutions. Maximum likelihood uses information on errors on one or both variables to choose a result between the two results of linear least-squares (direct and inverse fits). As a result, maximum likelihood cannot differ by more than one linear-bisector sigma from ordinary least-squares and even less from a linear-bisector fit. As we have shown that this linear-bisector sigma is about 1/3 of the Bayesian sigma, the maximum likelihood results should also be about 1/3 of the Bayesian values. Our results for the linear-bisector solutions thus suggest that the maximum likelihood method would yield distances and radii in the infrared surface brightness method that are unbiased if the uncertainties in the data are well understood. Unfortunately, these uncertainties are often not well understood. First, Barnes et al. (2003) showed that quoted uncertainties in Cepheid photometry and radial velocities are usually underestimated. Second, maximum-likelihood calculations usually adopt an approximation for the uncertainty in the displacements advocated by Balona (1977) for equally spaced velocity data. Our work demonstrates that this approximation does not apply to typical unequally-spaced radial velocity curves; if it did apply, the linear-bisector method would have yielded uncertainties in distance and radius close to those of the Bayesian MCMC calculation. Thus we expect the maximum likelihood uncertainties to be underestimated as are the linear-bisector uncertainties and for the same reasons.
TGB and WHJ gratefully acknowledge financial support for this work from McDonald Observatory and the Department of Astronomy of the University of Texas at Austin. WPG acknowledges financial support for this work from the Chilean Center for Astrophysics FONDAP 15010003. |
warning/0506/math0506556.html | ar5iv | text | # Numerical Solutions of a Boundary Value Problem for the Anomalous Diffusion Equation with the Riesz Fractional Derivative
## 1 Introduction
Anomalous diffusion is a phenomenon strongly connected with the interactions within complex and non-homogeneous background. This phenomenon is observed in transport of fluid in porous materials, in the chaotic heat baths, amorphous semiconductors, particle dynamics inside polymer network, two-dimensional rotating flow and also in econophysics. Phenomenon of anomalous diffusion deviates from the standard diffusion behaviour. In opposite to standard diffusion where linear form in the mean square displacement $`x^2\left(t\right)k_1t`$ of diffusing particle over time occurs, anomalous diffusion is characterized by the non-linear one $`x^2\left(t\right)k_\gamma t^\gamma `$, for $`\gamma (0,2]`$. In this phenomenon may exist dependence $`x^2\left(t\right)\mathrm{}`$ , which is characterized by occurrence of rare but extremely large jumps of diffusing particle – well-known as the Levy motion or the Levy flights. Ordinary diffusion follows Gaussian statistics and Fick’s second law for finding running process at time $`t`$ whereas anomalous diffusion follows non-Gaussian statistic or can be interpreted as the Levy stable densities.
Many authors proposed models which base on linear and non-linear forms of differential equations. Such models can simulate anomalous diffusion but they don’t reflect its real behaviour. Several authors apply fractional calculus in modelling of this type of diffusion. This means that time and spatial derivatives in the classical diffusion equation are replaced by fractional ones. In comparison to derivatives of integer order, which depend on the local behaviour of the function, derivatives of fractional order accumulate the whole history of this function.
## 2 Mathematical background
In this paper, we consider an equation in the following form
$$\frac{}{t}C(x,t)=k_\alpha \frac{^\alpha }{\left|x\right|^\alpha }C(x,t)\text{}t0\text{}x\text{}$$
(1)
where $`C(x,t)`$ is a field variable, $`\frac{^\alpha }{\left|x\right|^\alpha }C(x,t)`$ is the Riesz-Feller fractional operator , $`\alpha `$ is the real order of this operator, $`k_\alpha `$ is the coefficient of generalized (anomalous) diffusion with the unit of measure $`\left[m^\alpha /s\right]`$. According to the Riesz-Feller fractional operator for $`0<\alpha 2`$, $`\alpha 1`$ for one-variable function $`u(x)`$ is
$$\begin{array}{cc}\hfill \frac{^\alpha }{\left|x\right|^\alpha }u(x)={}_{x}{}^{}D_{\theta }^{\alpha }u\left(x\right)=& [c_L(\alpha ,\theta )_{\mathrm{}}D_x^\alpha u\left(x\right)\hfill \\ & +c_R(\alpha ,\theta )_xD_+\mathrm{}^\alpha u\left(x\right)]\text{,}\hfill \end{array}$$
(2)
where
$${}_{\mathrm{}}{}^{}D_{x}^{\alpha }u\left(x\right)=\{\begin{array}{cc}\frac{d}{dx}\left[{}_{\mathrm{}}{}^{}I_{x}^{1\alpha }u\left(x\right)\right],\hfill & \text{for}0<\alpha 1\text{,}\hfill \\ \frac{d^2}{dx^2}\left[{}_{\mathrm{}}{}^{}I_{x}^{2\alpha }u\left(x\right)\right],\hfill & \text{for}1<\alpha 2\text{,}\hfill \end{array}$$
(3)
$${}_{x}{}^{}D_{+\mathrm{}}^{\alpha }u\left(x\right)=\{\begin{array}{cc}\frac{d}{dx}\left[{}_{x}{}^{}I_{+\mathrm{}}^{1\alpha }u\left(x\right)\right],\hfill & \text{for}0<\alpha 1\text{,}\hfill \\ \frac{d^2}{dx^2}\left[{}_{x}{}^{}I_{+\mathrm{}}^{2\alpha }u\left(x\right)\right],\hfill & \text{for}1<\alpha 2\text{.}\hfill \end{array}$$
(4)
and coefficients $`c_L(\alpha ,\theta )`$, $`c_R(\alpha ,\theta )`$ (for $`0<\alpha 2`$, $`\alpha 1`$, and for $`\left|\theta \right|\mathrm{min}(\alpha ,2\alpha )`$), are defined as
$$c_L(\alpha ,\theta )=\frac{\mathrm{sin}{\displaystyle \frac{\left(\alpha \theta \right)\pi }{2}}}{\mathrm{sin}\left(\alpha \pi \right)}\text{}c_R(\alpha ,\theta )=\frac{\mathrm{sin}{\displaystyle \frac{\left(\alpha +\theta \right)\pi }{2}}}{\mathrm{sin}\left(\alpha \pi \right)}\text{.}$$
(5)
The fractional operators of order $`\alpha `$: $`{}_{\mathrm{}}{}^{}I_{x}^{\alpha }u\left(x\right)`$ and $`{}_{x}{}^{}I_{\mathrm{}}^{\alpha }u\left(x\right)`$ are defined as the left- and right-side of Weyl fractional integrals which definitions are
$${}_{\mathrm{}}{}^{}I_{x}^{\alpha }u\left(x\right)=\frac{1}{\mathrm{\Gamma }(\alpha )}_{\mathrm{}}^x\frac{u\left(\xi \right)}{(x\xi )^{1\alpha }}𝑑\xi \text{,}$$
(6)
$${}_{x}{}^{}I_{\mathrm{}}^{\alpha }u\left(x\right)=\frac{1}{\mathrm{\Gamma }(\alpha )}_x^{\mathrm{}}\frac{u\left(\xi \right)}{(\xi x)^{1\alpha }}𝑑\xi \text{.}$$
(7)
Considering Eqn (1) we obtain the classical diffusion equation for $`\alpha =2`$, i.e. the heat transfer equation. If $`\alpha =1`$, and the parameter of skewness $`\theta `$ admits extreme values in (5), the transport equation is noted. Therefore we assume variations of the parameter $`\alpha `$ within the range $`0<\alpha 2`$. Analysing behaviour of the parameter $`\alpha <2`$ in Eqn (1), we found some combination between transport and propagation processes.
For analytic solution of Eqn (1) we can apply Green functions . We numerically solve Eqn (1) when additional non-linear term may occur. Some numerical methods used in solution of fractional differential equations can be found in . However they apply the infinite domain.
In this work we will consider Eqn (1) limited for $`1<\alpha 2`$ in one dimensional domain $`\mathrm{\Omega }:LxR`$ with the boundary-value conditions of the first kind (the Dirichlet conditions) as
$$\{\begin{array}{cc}x=L:\hfill & C(L,t)=g_L\left(t\right),\hfill \\ x=R:\hfill & C(R,t)=g_R\left(t\right),\hfill \end{array}t>0,$$
(8)
and with the initial-value condition
$$C(x,t)|_{t=0}=c_0\left(x\right)\text{.}$$
(9)
## 3 Numerical method
According to the finite difference method we consider a discrete from of Eqn (1) both in time and space. In the previous work we solved numerically the anomalous diffusion equation similar to the Eqn (1) with the time-fractional derivative. We called this method FFDM (Fractional FDM). The problem of solving of Eqn (1) lies in properly approximation of the Riesz-Feller derivative (2) in numerical scheme.
### 3.1 Approximation of the Riesz-Feller derivative
We begin numerical analysis from discrete forms of operators (6) and (7). We introduce homogenous spatial grid $`\mathrm{}<\mathrm{}<x_{i2}<x_{i1}<x_i<x_{i+1}<x_{i+2}<\mathrm{}<\mathrm{}`$ with the step $`h=x_kx_{k1}`$ and we denote value of function $`u`$ in the point $`x_k`$ as $`u_k=u\left(x_k\right)`$, for $`k`$. In order to simplify notations we take here the function of one variable. For numerical integration scheme we assumed the trapezoidal rule. The integral (6) in point $`x_i`$ of the grid is replaced by the sum of discrete integrals as
$${}_{\mathrm{}}{}^{}I_{x_i}^{\alpha }u_i=\frac{1}{\mathrm{\Gamma }(\alpha )}\underset{k=0}{\overset{\mathrm{}}{}}\underset{x_{ik1}}{\overset{x_{ik}}{}}\frac{u\left(\xi \right)}{(x_i\xi )^{1\alpha }}𝑑\xi \text{,}$$
(10)
and using linear interpolation of function $`u`$ in every sub-interval $`[x_{ik1},x_{ik}]`$
$$u^{}\left(\xi \right)=\frac{u_{ik}u_{ik1}}{h}\xi +\frac{u_{ik1}x_{ik}u_{ik}x_{ik1}}{h}$$
(11)
we have
$`{}_{\mathrm{}}{}^{}I_{x_i}^{\alpha }u_i`$ $``$ $`{\displaystyle \frac{1}{\mathrm{\Gamma }(\alpha )}}{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{x_{ik1}}{\overset{x_{ik}}{}}}{\displaystyle \frac{u^{}\left(\xi \right)}{(x_i\xi )^{1\alpha }}}𝑑\xi `$
$`=`$ $`{\displaystyle \frac{1}{\mathrm{\Gamma }(\alpha )}}{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}[(u_{ik}u_{ik1})a_k^{\left(\alpha \right)}`$
$`+(u_{ik1}x_{ik}u_{ik}x_{ik1})b_k^{\left(\alpha \right)}]`$
where
$`a_k^{\left(\alpha \right)}`$ $`=`$ $`h^{\alpha 1}x_i{\displaystyle \frac{\left(k+1\right)^\alpha k^\alpha }{\alpha }}h^\alpha {\displaystyle \frac{\left(k+1\right)^{1+\alpha }k^{1+\alpha }}{1+\alpha }}\text{,}`$ (13)
$`b_k^{\left(\alpha \right)}`$ $`=`$ $`h^{\alpha 1}{\displaystyle \frac{\left(k+1\right)^\alpha k^\alpha }{\alpha }}\text{.}`$ (14)
After next transforms we can write
$${}_{\mathrm{}}{}^{}I_{x_i}^{\alpha }u_ih^\alpha \underset{k=0}{\overset{\mathrm{}}{}}u_{ik}v_k^{\left(\alpha \right)}$$
(15)
where
$`v_k^{\left(\alpha \right)}={\displaystyle \frac{1}{\mathrm{\Gamma }(2+\alpha )}}\times `$ (16)
$`\{\begin{array}{cc}1\text{,}\hfill & \text{ for }k=0\text{,}\hfill \\ \left(k+1\right)^{1+\alpha }2k^{1+\alpha }+\left(k1\right)^{1+\alpha }\text{,}\hfill & \text{ for }k=1,\mathrm{},\mathrm{}\text{.}\hfill \end{array}`$ (19)
Similar to previous considerations we approximate operator $`{}_{x}{}^{}I_{\mathrm{}}^{\alpha }u\left(x\right)`$ in the point $`x_i`$ and finally we obtain
$${}_{x_i}{}^{}I_{\mathrm{}}^{\alpha }u_ih^\alpha \underset{k=0}{\overset{\mathrm{}}{}}u_{i+k}v_k^{\left(\alpha \right)}\text{,}$$
(20)
where coefficients $`v_k^{\left(\alpha \right)}`$ have identical forms as (16).
In the next step we analyse operator (2). It can be expressed in the form (in order to simplify this we denote $`c_L=c_L(\alpha ,\theta )`$ and $`c_R=c_R(\alpha ,\theta )`$ )
$`{}_{x}{}^{}D_{\theta }^{\alpha }u\left(x\right)=`$ (21)
$`\left[c_L{\displaystyle \frac{d^2}{dx^2}}\left[{}_{\mathrm{}}{}^{}I_{x}^{2\alpha }u\left(x\right)\right]+c_R{\displaystyle \frac{d^2}{dx^2}}\left[{}_{x}{}^{}I_{+\mathrm{}}^{2\alpha }u\left(x\right)\right]\right]\text{.}`$
We used the central difference scheme for the second spatial derivative in the point $`x_i`$ and we obtain
$`{}_{x_i}{}^{}D_{\theta }^{\alpha }u_i`$ (22)
$`\begin{array}{c}[c_L{\displaystyle \frac{{}_{\mathrm{}}{}^{}I_{x}^{2\alpha }u_{i1}2_{\mathrm{}}I_x^{2\alpha }u_i+_{\mathrm{}}I_x^{2\alpha }u_{i+1}}{h^2}}\\ +c_R{\displaystyle \frac{{}_{x}{}^{}I_{+\mathrm{}}^{2\alpha }u_{i1}2_xI_+\mathrm{}^{2\alpha }u_i+_xI_+\mathrm{}^{2\alpha }u_{i+1}}{h^2}}]\text{.}\end{array}`$ (25)
After numerous transforms we obtain the final form as
$${}_{x_i}{}^{}D_{\theta }^{\alpha }u_i\frac{1}{h^\alpha }\underset{k=\mathrm{}}{\overset{\mathrm{}}{}}u_{i+k}w_k^{\left(\alpha \right)}\text{,}$$
(26)
where coefficients $`w_k^{\left(\alpha \right)}`$ are
$`w_k^{\left(\alpha \right)}={\displaystyle \frac{1}{\mathrm{\Gamma }\left(4\alpha \right)}}\times `$ (27)
$`\{\begin{array}{cc}(\left(\right|k|+2)^{3\alpha }4\left(\right|k|+1)^{3\alpha }+6|k|^{3\alpha }\hfill & \\ 4\left(\right|k|1)^{3\alpha }+\left(\right|k|2)^{3\alpha })c_L,\hfill & \text{for }k2\hfill \\ \left(3^{3\alpha }2^{5\alpha }+6\right)c_L+c_R,\hfill & \text{for }k=1\hfill \\ \left(2^{3\alpha }4\right)\left(c_L+c_R\right),\hfill & \text{for }k=0\hfill \\ \left(3^{3\alpha }2^{5\alpha }+6\right)c_R+c_L,\hfill & \text{for }k=1\hfill \\ ((k+2)^{3\alpha }4(k+1)^{3\alpha }+6k^{3\alpha }\hfill & \\ 4(k1)^{3\alpha }+(k2)^{3\alpha })c_R,\hfill & \text{for }k2\hfill \end{array}\text{.}`$ (35)
Assuming $`\alpha =2`$ and $`\theta =0`$ we have $`c_L(2,0)=`$ $`c_R(2,0)=\frac{1}{2}`$ and we obtain
$$w_k^{\left(2\right)}=\{\begin{array}{cc}0,\hfill & \text{for }k2\hfill \\ 1,\hfill & \text{for }k=1\hfill \\ 2,\hfill & \text{for }k=0\hfill \\ 1,\hfill & \text{for }k=1\hfill \\ 0,\hfill & \text{for }k2\hfill \end{array}\text{.}$$
(36)
These coeeficients are identical as for wide known the central difference scheme for the second derivative. Also when $`\alpha 1^+`$ and $`\theta =0`$ after arduous calculations of limits we obtain coefficients
$`w_k^{\left(1^+\right)}={\displaystyle \frac{1}{2\pi }}\times `$ (37)
$`\{\begin{array}{cc}\mathrm{ln}{\displaystyle \frac{\left(\left|k\right|+1\right)^{4\left(\left|k\right|+1\right)^2}\left(\left|k\right|1\right)^{4\left(\left|k\right|1\right)^2}}{\left(\left|k\right|+2\right)^{\left(\left|k\right|+2\right)^2}\left|k\right|^{6k^2}\left(\left|k\right|2\right)^{\left(\left|k\right|2\right)^2}}},\hfill & \text{ for }k2,\hfill \\ 16\mathrm{ln}29\mathrm{ln}3,\hfill & \text{ for }k=1,\hfill \\ 8\mathrm{ln}2,\hfill & \text{ for }k=0,\hfill \\ 16\mathrm{ln}29\mathrm{ln}3,\hfill & \text{ for }k=1,\hfill \\ \mathrm{ln}{\displaystyle \frac{\left(k+1\right)^{4\left(k+1\right)^2}\left(k1\right)^{4\left(k1\right)^2}}{\left(k+2\right)^{\left(k+2\right)^2}k^{6k^2}\left(k2\right)^{\left(k2\right)^2}}},\hfill & \text{ for }k2.\hfill \end{array}`$ (43)
In literature didn’t find exact values of approximating coefficients. When $`\alpha =1`$ the Riesz-Feller operator is singular, hence the problem. Numerous works of Gorenflo and Mainardi i.e. propose various ways which determine values of the coefficients $`w_k^{\left(\alpha \right)}`$ (i.e. based on the Grünwald-Letnikov discretization) but they don’t provide continuity in the interval $`\alpha (1,2]`$. The coefficients (23) can approximate the Cauchy process when we use (23) in numerical calculations.
### 3.2 Fractional FDM
While discretization of the Riesz-Feller derivative in space is done, in this subsection we describe the finite difference method for the equation of anomalous diffusion (1). Here we restrict this solution to one dimensional space. In comparison with the standard diffusion equation where discretization of the second derivative in space can be approximated by the central difference of second order, we will use generalized scheme given by formula (20). The differences appear in setting of boundary conditions.
We shall introduce a temporal grid $`0=t^0<t^1<\mathrm{}<t^f<t^{f+1}<\mathrm{}<`$$`\mathrm{}`$ with the step $`\mathrm{\Delta }t=t^{f+1}t^f`$ and we denote value of the function $`C(x,t)`$ in the point $`x_k`$ at the moment of time $`t^f`$ as $`C_k^f=C(x_k,t^f)`$ for $`k`$ and $`f`$.
#### 3.2.1 Pure initial value problem
In the explicit scheme of the FDM we replaced Eqn (1) by the following formula
$$\frac{C_i^{f+1}C_i^f}{\mathrm{\Delta }t}=K_\alpha \frac{1}{h^\alpha }\underset{k=\mathrm{}}{\overset{\mathrm{}}{}}C_{i+k}^fw_k^{\left(\alpha \right)}\text{.}$$
(44)
After simplification finally we obtained
$$C_i^{f+1}=\underset{k=\mathrm{}}{\overset{\mathrm{}}{}}C_{i+k}^fp_k^{\left(\alpha \right)}\text{,}$$
(45)
where coefficients $`p_k^{\left(\alpha \right)}`$ are
$$p_k^{\left(\alpha \right)}=\{\begin{array}{cc}1+K_\alpha \frac{\mathrm{\Delta }t}{h^\alpha }w_0^{\left(\alpha \right)},\hfill & \text{for }k=0\text{,}\hfill \\ K_\alpha \frac{\mathrm{\Delta }t}{h^\alpha }w_k^{\left(\alpha \right)},\hfill & \text{for }k0.\hfill \end{array}$$
(46)
Using simple calculations one may proof, that arise the following relationship
$$\underset{k=\mathrm{}}{\overset{\mathrm{}}{}}p_k^{\left(\alpha \right)}=1\text{.}$$
(47)
In order to determine stability of the explicit scheme the coefficient (26) for $`k=0`$ in formula (25) should be positive
$$p_0^{\left(\alpha \right)}=1+K_\alpha \frac{\mathrm{\Delta }t}{h^\alpha }w_0^{\left(\alpha \right)}>0\text{.}$$
(48)
Hence we fixed the maximum length of the step $`\mathrm{\Delta }t`$ as
$$\mathrm{\Delta }t<\frac{h^\alpha }{K_\alpha w_0^{\left(\alpha \right)}}=\frac{h^\alpha \mathrm{\Gamma }\left(4\alpha \right)}{K_\alpha \left(2^{3\alpha }4\right)\left(c_L(\alpha ,\theta )+c_R(\alpha ,\theta )\right)}\text{.}$$
(49)
The initial condition (9) is introduced directly to every grid nodes at the first step $`t=t^0`$. This determines values of the function $`C`$ as
$$C_i^0=c_0\left(x_i\right)\text{.}$$
(50)
In unbounded domains the implicit method isn’t easily applicable because it generates infinite dimensions of all matrices. Thus one usually seeks improved difference equations within the explicit scheme.
#### 3.2.2 Boundary-initial value problem
Presenting numerical solution (25) with included unbounded domain $`\mathrm{}<x<\mathrm{}`$ has no practical implementations in computer simulations.
Now, we present solution of this problem on the finite domain $`\mathrm{\Omega }:LxR`$ with boundary conditions (8). We divide this domain $`\mathrm{\Omega }`$ into $`N`$ sub-domains with $`h=(RL)/N`$. Figure 1 shows modified spatial grid.
Here we can observe additional ’virtual’ points in the grid placed outside of the domain $`\mathrm{\Omega }`$. In order to introduce the Dirichlet boundary conditions we proposed treatment which bases on assumption that values of the function $`C`$ in outside points are identical as values in the boundary nodes $`x_0`$ or $`x_N`$
$$C(x_k,t)=\{\begin{array}{cc}C(x_0,t)=g_L\left(t\right)\text{,}\hfill & \text{for }k<0\text{,}\hfill \\ C(x_N,t)=g_R\left(t\right)\text{,}\hfill & \text{for }k>N\text{.}\hfill \end{array}$$
(51)
On the base of previous considerations we modify expression (20) for discretization of the Riesz-Feller derivative. Thus we have
$`{}_{x_i}{}^{}D_{\theta }^{\alpha }C(x_i,t)`$ $``$ $`{\displaystyle \frac{1}{h^\alpha }}[{\displaystyle \underset{k=i}{\overset{Ni}{}}}C(x_{i+k},t)w_k^{\left(\alpha \right)}`$ (52)
$`+g_L\left(t\right)s_{L}^{}{}_{i}{}^{\left(\alpha \right)}+g_R\left(t\right)s_{R}^{}{}_{i}{}^{\left(\alpha \right)}],`$
for $`i=1,\mathrm{},N1`$, where
$`s_{L}^{}{}_{i}{}^{\left(\alpha \right)}={\displaystyle \underset{k=\mathrm{}}{\overset{i1}{}}}w_k^{\left(\alpha \right)}={\displaystyle \frac{1}{\mathrm{\Gamma }\left(4\alpha \right)}}\times `$ (53)
$`\left[\left(i+2\right)^{3\alpha }+3\left(i+1\right)^{3\alpha }3i^{3\alpha }+\left(i1\right)^{3\alpha }\right]c_L,`$
$`s_{R}^{}{}_{i}{}^{\left(\alpha \right)}={\displaystyle \underset{k=Ni+1}{\overset{\mathrm{}}{}}}w_k^{\left(\alpha \right)}={\displaystyle \frac{1}{\mathrm{\Gamma }\left(4\alpha \right)}}[(Ni+2)^{3\alpha }`$ (54)
$`+3(Ni+1)^{3\alpha }3(Ni)^{3\alpha }+(Ni1)^{3\alpha }]c_R.`$
Putting this expression to Eqn (1) we obtain a finite difference scheme depending on weighting factor $`\sigma `$. Here we assumed
$`g_L^{f+\frac{1}{2}}`$ $`=`$ $`g_L\left(t^{f+\frac{1}{2}}\right)=g_L\left(\mathrm{\Delta }t\left(f+{\displaystyle \frac{1}{2}}\right)\right),`$ (55)
$`g_R^{f+\frac{1}{2}}`$ $`=`$ $`g_R\left(t^{f+\frac{1}{2}}\right)=g_R\left(\mathrm{\Delta }t\left(f+{\displaystyle \frac{1}{2}}\right)\right)`$ (56)
in order to simplify the numerical scheme. For internal nodes $`x_i`$, $`i=1,\mathrm{},N1`$ we have
$`{\displaystyle \frac{C_i^{f+1}C_i^f}{\mathrm{\Delta }t}}`$ $`=`$ $`K_\alpha {\displaystyle \frac{1}{h^\alpha }}[{\displaystyle \underset{k=i}{\overset{Ni}{}}}(\sigma C_{i+k}^f+(1\sigma )C_{i+k}^{f+1})w_k^{\left(\alpha \right)}`$ (57)
$`+g_L^{f+\frac{1}{2}}s_L{}_{i}{}^{\left(\alpha \right)}+g_R^{f+\frac{1}{2}}s_R{}_{i}{}^{\left(\alpha \right)}],`$
and for the boundary nodes $`x_0`$ and $`x_N`$:
$`C_0^{f+1}`$ $`=`$ $`g_L^{f+\frac{1}{2}},`$ (58)
$`C_N^{f+1}`$ $`=`$ $`g_R^{f+\frac{1}{2}}.`$ (59)
The method is explicit for $`\sigma =1`$ and partially implicit for $`0<\sigma <1`$ and with $`\sigma =0`$ being fully implicit. In literature this method is known as the $`\sigma `$-method for parabolic equations.
Above scheme described by expressions (37)-(39) can be written in matrix form as
$$𝐀𝐂^{𝐟+\mathrm{𝟏}}=𝐁\text{,}$$
(60)
where
$$𝐀=\left[\begin{array}{cccccccc}1& 0& 0& 0& \mathrm{}& 0& 0& 0\\ a_1& 1+a_0& a_1& a_2& \mathrm{}& a_{N3}& a_{N2}& a_{N1}\\ a_2& a_1& 1+a_0& a_1& \mathrm{}& a_{N4}& a_{N3}& a_{N2}\\ a_3& a_2& a_1& 1+a_0& \mathrm{}& a_{N5}& a_{N4}& a_{N2}\\ a_4& a_3& a_2& a_1& \mathrm{}& a_{N6}& a_{N3}& a_{N4}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ a_{N+2}& a_{N+3}& a_{N+4}& a_{N+5}& \mathrm{}& 1+a_0& a_1& a_2\\ a_{N+1}& a_{N+2}& a_{N+3}& a_{N+4}& \mathrm{}& a_1& 1+a_0& a_1\\ 0& 0& 0& 0& \mathrm{}& 0& 0& 1\end{array}\right],$$
(61)
$$𝐁=\left[\begin{array}{c}g_L^{f+\frac{1}{2}}\\ b_1\\ b_2\\ b_3\\ b_4\\ \mathrm{}\\ b_{N2}\\ b_{N1}\\ g_R^{f+\frac{1}{2}}\end{array}\right],$$
(62)
with
$`a_j`$ $`=`$ $`\left(\sigma 1\right)K_\alpha {\displaystyle \frac{\mathrm{\Delta }t}{h^\alpha }}w_j^{\left(\alpha \right)},\text{ for }j=N+1,\mathrm{},N1\text{,}`$ (63)
$`b_j`$ $`=`$ $`C_j^f+K_\alpha {\displaystyle \frac{\mathrm{\Delta }t}{h^\alpha }}[g_L^{f+\frac{1}{2}}s_{L}^{}{}_{j}{}^{\left(\alpha \right)}+g_R^{f+\frac{1}{2}}s_{R}^{}{}_{j}{}^{\left(\alpha \right)}`$ (64)
$`+\sigma {\displaystyle \underset{k=j}{\overset{Nj}{}}}C_{i+k}^fw_k^{\left(\alpha \right)}],\text{ for }j=1,\mathrm{},N1\text{.}`$
and $`𝐂^{𝐟+\mathrm{𝟏}}`$ is the vector of unknown function’s values $`C`$ at the time $`t^{f+1}`$.
Particular case of above scheme (37) is the explicit scheme (for $`\sigma =1`$) which may be simplified to
$$C_i^{f+1}=\{\begin{array}{cc}g_L^{f+\frac{1}{2}},\hfill & \text{for }i=0,\hfill \\ K_\alpha \frac{\mathrm{\Delta }t}{h^\alpha }\left(g_L^{f+\frac{1}{2}}s_{L}^{}{}_{i}{}^{\left(\alpha \right)}+g_R^{f+\frac{1}{2}}s_{R}^{}{}_{i}{}^{\left(\alpha \right)}\right)\hfill & \\ +\underset{k=i}{\overset{ni}{}}C_{i+k}^fp_k^{\left(\alpha \right)}\text{}\hfill & \text{for }i=1,\mathrm{},N1,\hfill \\ g_R^{f+\frac{1}{2}},\hfill & \text{for }i=N,\hfill \end{array}$$
(65)
with $`p_k^{\left(\alpha \right)}`$ defined by formula (26).
We can observe that boundary conditions influence to all values of the function in every node. In opposite to the second derivative over space which is approximated locally, the characteristic feature of Riesz-Feller and other fractional derivatives is dependence on values of all domain points. For $`\alpha =2`$ and $`\theta =0`$ our scheme is identically as wide known and used the forward difference in time and central difference in space scheme (FTCS) .
The skewness parameter $`\theta `$ has great significance influence on the solution. For $`\alpha 1^+`$ and $`\theta \pm 1^+`$ one can obtain the classical hyperbolic equation, i.e. the first order wave equation (the transport equation). In this case our scheme tends to the known Euler’s forward time and central space (FTCS) approximation of Eqn (1). Unfortunately this is unconditionally unstable and therefore this is disadvantage this method.
Proposed numerical scheme makes a bridge between Gaussian and Cauchy processes. Our scheme is also a bridge between diffusion and transport phenomena.
## 4 Simulation results
In this section we present results of calculation. In all presented simulations we assumed $`k_\alpha =1m^\alpha /s`$ and the length of 1D domain $`l=1m`$. Figure 1 shows two charts over space (one in the logarithmic scale) with absorbing boundary $`C(x,t)|_{x=0}=C(x,t)|_{x=1}=0`$. On these plots solutions for different values of parameter $`\alpha \mathrm{\hspace{0.17em}1.01},1.5,2`$ at time $`t=0,0.01,0.3s`$ for $`\theta =0`$ are presented.
Figure 2 presents another example of the solution which differs from example presented by the Fig. 1 (boundary conditions $`C(x,t)|_{x=0}=C(x,t)|_{x=1}=100`$ and initial condition $`C(x,t)|_{t=0}=0`$ ). In both cases we observe diffusion process arising in different way. The last example reflects case when the parameter of skewness is $`\theta =0.5`$ and $`\alpha =1.4`$. Figure 3 shows a diffusion transport process over space at different moments of time.
## 5 Conclusions
In summary, we proposed the fractional finite difference method for fractional diffusion equation with the Riesz-Feller fractional derivative which is extension to the standard diffusion. We analysed a linear case of diffusion equation and in the future we will work on non-linear cases. We obtained the implicit and explicit FDM schemes which generalise classical schemes of FDM for the diffusion equation. Moreover, for $`\alpha =2`$ our solution equals to the classical finite difference method.
Analysing plots included in this work, we can see that in the case $`\alpha <2`$ (the Levy flight) diffusion is slower then the standard diffusion (Brownian motion) in the initial time. Nevertheless, when we analyse the probability density function we observe a long tail of distribution in the long time limit. In this way we can simulate same rare and extreme events which are characterised by arbitrary very large values of particle jumps.
Analysing changes in the skewness parameter $`\theta `$ we observed interesting behaviour in solution. For $`\alpha 1^+`$ and for $`\theta \pm 1^+`$ we obtained the first order wave equation. For $`\theta (0,1)`$ (with restrictions to order $`\alpha `$) we generate a class of non-symmetric probability density functions. |
warning/0506/astro-ph0506451.html | ar5iv | text | # TOWARDS AN EMPIRICAL DETERMINATION OF THE ZZ CETI INSTABILITY STRIP
## 1 INTRODUCTION
The ZZ Ceti stars represent a class of variable white dwarfs whose optical spectra are dominated by hydrogen lines (DA stars). They occupy a narrow region in the $`T_{\mathrm{eff}}`$-$`\mathrm{log}g`$ plane known as the ZZ Ceti instability strip, with an average effective temperature around $`T_{\mathrm{eff}}11,600`$ K and a width of roughly 1000 K. A precise determination of the hot and cool boundaries of this instability strip may eventually provide important constraints about the structure of the outer layers of DA white dwarfs. For instance, it has been originally shown by Winget et al. (1982) that the location of the blue edge is sensitive to the convective efficiency in the hydrogen zone, which led Fontaine, Tassoul, & Wesemael (1984) to propose using this property as a potential calibrator of the mixing-length theory in pulsating white dwarfs. Similarly, the location of the red edge may help us understand the mechanism responsible for the disappearance of the ZZ Ceti phenomenon at low temperatures, which seems to be related to either convective mixing of the hydrogen outer layer with the deep helium envelope or the interaction of pulsation with convection (Tassoul et al., 1990). Also of utmost importance is to determine whether all white dwarfs within the ZZ Ceti instability strip are pulsators. If the strip is indeed pure, as first suggested by Fontaine et al. (1982), ZZ Ceti stars would necessarily represent a phase through which all DA stars must evolve, and thus the results from asteroseismological studies might provide constraints on the properties not only of known ZZ Ceti stars, but on the whole population of DA stars as well.
Determinations of the boundaries of the ZZ Ceti instability strip prior to 1991 have been nicely summarized by Wesemael et al. (1991) who discuss the results from various observational techniques, both photometric and spectroscopic. Among the first photometric studies were those conducted using Strömgren photometry by McGraw (1979) and later by Fontaine et al. (1985). Both analyses made it evident that ZZ Ceti stars formed a rather homogeneous class of DA white dwarfs in color-color diagrams, a result that was not obvious from prior analyses based on broad-band colors. Multichannel spectrophotometric data of ZZ Ceti stars obtained by Greenstein (1976) have been analyzed by Fontaine et al. (1982), Greenstein (1982), and by Weidemann & Koester (1984) using slightly different absolute flux calibrations.
Later on, Wesemael et al. (1986), Lamontagne et al. (1987), and Lamontagne et al. (1989) have used ultraviolet spectra obtained by the IUE satellite as an independent method of measuring the effective temperature of ZZ Ceti stars. In their analysis they assumed a value of $`\mathrm{log}g=8`$ for each star but also mentioned that this assumption could be a source of uncertainty as several ZZ Ceti stars showed signs of having $`\mathrm{log}g`$ significantly higher or lower (e.g., G226-29 and Ross 548, respectively). Finally, Daou et al. (1990) have carried out the first analysis of a set of ZZ Ceti stars using the spectroscopic technique where optical spectroscopic observations of the individual Balmer lines are fitted with synthetic spectra to obtain measures of both $`T_{\mathrm{eff}}`$ and $`\mathrm{log}g`$.
The effective temperatures for the ZZ Ceti stars inferred from these photometric and spectroscopic studies are in fairly good agreement according to Figure 1 of Wesemael et al. (1991), with the blue edge in the range $`T_{\mathrm{eff}}=12,13013,500`$ K and the red edge in the range $`T_{\mathrm{eff}}=10,00011,740`$ K. However, this apparent agreement has been seriously questioned by Bergeron et al. (1992b) who examined the effects of different convective efficiencies on the optical spectra of DA white dwarfs in the vicinity of the ZZ Ceti instability strip. The results of their calculations showed that the predicted absolute fluxes, color indices, and equivalent widths are sensitive to the convective efficiency in the range $`T_{\mathrm{eff}}800015,000`$ K, with a maximum sensitivity around 13,000 K. Hence, without a detailed knowledge of the convective efficiency in the atmosphere of ZZ Ceti stars, the results from all previous photometric and spectroscopic analyses had to be considered uncertain.
This problem of the convective efficiency in the atmosphere of ZZ Ceti stars has been tackled by Bergeron et al. (1995c, B95 hereafter) who used optical spectroscopic observations combined with UV energy distributions to show that the so-called ML2/$`\alpha =0.6`$ parametrization of the mixing-length theory provides the best internal consistency between optical and UV effective temperatures, trigonometric parallaxes, $`V`$ magnitudes, and gravitational redshifts. With the atmospheric convective efficiency properly parameterized, the spectroscopic technique could now yield atmospheric parameters $`T_{\mathrm{eff}}`$ and $`\mathrm{log}g`$ for the ZZ Ceti stars that were not only accurate in a relative sense, but in an absolute sense as well. Hence it was possible for the first time to demonstrate that the boundaries of the ZZ Ceti instability strip were a function of both the effective temperature and the surface gravity of the star. Our knowledge of the boundaries of the ZZ Ceti instability strip prior to the study of Mukadam et al. (2004b) discussed below is summarized in Figure 4 of Bergeron et al. (2004). The ZZ Ceti stars occupy a trapezoidal region in the $`T_{\mathrm{eff}}\mathrm{log}g`$ plane, with the blue edge showing a stronger dependence on the surface gravity than the red edge does. Consequently the width of the instability strip is also gravity-dependent, with $`\mathrm{\Delta }T_{\mathrm{eff}}800`$ K at $`\mathrm{log}g=7.5`$ and nearly twice as wide at $`\mathrm{log}g=8.5`$.
As mentioned above, the assessment of the purity of the instability strip is also of considerable interest. More than twenty years ago, Fontaine et al. (1982) have argued from their study of multichannel spectrophotometric data that the strip is most likely pure, and that ZZ Ceti stars therefore represent an evolutionary phase through which all DA white dwarfs must pass. This conclusion is strongly supported by our spectroscopic analysis of the 36 known ZZ Ceti stars shown in Figure 4 of Bergeron et al. (2004). The latter also included 54 known nonvariable white dwarfs that were all found to lie clearly outside the empirical instability strip. We note that, prior to this effort, the purity of the instability strip had been questioned repeatedly (Dolez et al., 1991; Kepler & Nelan, 1993; Kepler et al., 1995; Silvotti et al., 1997; Giovannini et al., 1998).
More recently, Mukadam et al. (2004a) reported the discovery of 35 new ZZ Ceti stars from the Sloan Digital Sky Survey (SDSS) along with a large number of stars found to be photometrically constant. Mukadam et al. (2004b) used the results from this sample of both variable and nonvariable DA stars to “redefine” the location of the instability strip and to assess its purity. Although their determinations of $`T_{\mathrm{eff}}`$ and $`\mathrm{log}g`$ for these new ZZ Ceti stars place virtually all of the variables within the instability strip defined by Bergeron et al. (2004), with a possible offset due to the use of a different set of model spectra (see § 3), they also found a large fraction of nonvariable stars within the strip. These results are clearly at odds with the conclusions from our work during the last ten years.
In this respect, we have been gathering over the recent years optical spectroscopic observations for all known nonvariable DA white dwarfs with the goal of (1) constraining the location of the boundaries of the ZZ Ceti instability strip not only by analyzing the variable stars within the strip itself, but also the photometrically constant stars in its vicinity, and (2) increasing the statistical significance of the purity of the empirical instability strip. Some partial results from this endeavor have been reported in Bergeron et al. (2004). Here we present the results of our entire sample in § 2, which include the discovery of a new ZZ Ceti star. In § 3 we revisit the results of Mukadam et al. (2004b) for the variable and nonvariable DA stars uncovered in the SDSS. We then report in § 4 on preliminary results of a much broader spectroscopy survey of the white dwarf catalog of McCook & Sion (1999). Our conclusions follow in § 5.
## 2 PHOTOMETRIC SAMPLE
### 2.1 Spectroscopic Observations
Our sample of photometrically constant DA stars is composed of 121 objects gathered from various sources. Firstly, we have searched the literature for all mentions of DA white dwarfs observed in high-speed photometry and where no variations were detected. These include two Ph.D. theses (McGraw, 1977; Giovannini, 1996) and several studies of the instability strip including those of Dolez et al. (1991), Kepler et al. (1995), and Giovannini et al. (1998). Another source consists of previously unpublished data from various observing campaigns conducted over the years by two of us (GF and PB) and collaborators. We were also able to include 4 white dwarfs identified in the Hamburg Quasar Survey and reported to be constant by Mukadam et al. (2004a). Our sample does not include, however, stars whose nonvariability has recently come to our attention such as those reported by Silvotti et al. (2005) nor those discovered in the SDSS by Mukadam et al. (2004a) and Mullally et al. (2005).
Our sample of 121 nonvariable DA stars is listed in Table 1 in order of increasing right ascension. About 30% of the spectra in this sample were already available from the previous spectroscopic analyses of Bergeron et al. (1992a) and Bergeron et al. (1995a). These spectra had been secured using our standard setup at the Steward Observatory 2.3 m telescope equipped with the Boller & Chivens spectrograph. The $`4.^{\prime \prime }5`$ slit together with the 600 line mm<sup>-1</sup> grating blazed at 3568 Å in first order provides a spectral coverage from about 3000 to 5250 Å at a resolution of $``$ 6 Å FWHM. An additional 40 spectra were provided to us by C. Moran (1999, private communication); these have a comparable spectral coverage but at a slightly better resolution of $`3`$ Å FWHM. Seven spectra from the southern hemisphere are taken from the analyses of Bragaglia et al. (1995) and Bergeron et al. (2001). Finally, high signal-to-noise ratio (S/N) optical spectra for 36 objects were obtained specifically for the purpose of this project during four observing runs in 2003 and 2004, using again the Steward Observatory facility.
### 2.2 Fitting Procedure
The method used for fitting the spectroscopic observations relies on the so-called spectroscopic technique developed by Bergeron et al. (1992a), and which has been refined by B95 and more recently by Liebert et al. (2005, LBH hereafter). The most important improvement of the method is the way the continuum used to normalize individual Balmer lines is defined. The approach is slightly different depending on the temperature range in question. For stars in the interval 16,000 $`T_{\mathrm{eff}}`$ 9000 K, pseudo-Gaussian profiles are used whereas outside this temperature range synthetic spectra are utilized to determine the continuum (see Fig. 4 of LBH). Once the Balmer lines are normalized properly, we proceed to fit them with a grid of synthetic spectra derived from model atmospheres with a pure hydrogen composition. Our grid covers a range between $`T_{\mathrm{eff}}=1500`$ K and 140,000 K by steps of 500 K at low temperatures ($`T_{\mathrm{eff}}<17,000`$ K) and 5000 K at high temperatures ($`T_{\mathrm{eff}}>20,000`$ K), and a range in $`\mathrm{log}g`$ between 6.5 and 9.5 by steps of 0.5 dex (steps of 0.25 dex are used between 8000 K and 17,000 K where Balmer lines reach their maxima). For models where convective energy transport becomes important, we adopt the ML2/$`\alpha =0.6`$ parametrization of the mixing-length theory, as prescribed by B95.
One of the trickiest aspects of fitting optical spectra near the ZZ Ceti instability strip is the fact that we overlap the temperature interval over which the equivalent widths of the Balmer lines reach their maximum near $`T_{\mathrm{eff}}1314,000`$ K (see, e.g., Fig. 4 of B95). Hence, in some cases, the minimization procedure allows two acceptable solutions, one on each side of this maximum. When the true effective temperature of the star is more than $`2000`$ K away from the maximum, it is possible from a simple visual inspection of the fits to discriminate between the cool and the hot solutions. Indeed, for identical equivalent widths, the Balmer lines on the cool side of the maximum have deeper line cores. For stars in the range $`T_{\mathrm{eff}}11,50016,000`$ K we rely on the slopes of the observed and theoretical spectra normalized to unity at 4600 Å to discriminate between both solutions. As the slope of the energy distribution changes rapidly with temperature, it becomes relatively easy to decide which solution to adopt. Finally, whenever possible, our choice of solution has been confirmed by comparing multichannel, Strömgren, or Johnson photometry published in McCook & Sion (1999) with the theoretical color predictions of Bergeron et al. (1995b).
LBH used multiple spectroscopic observations of individual white dwarfs to estimate the external uncertainties of the fitted atmospheric parameters obtained from the spectroscopic technique (see their Fig. 8). Their estimate of the external error of each fitted parameter is 1.2% in $`T_{\mathrm{eff}}`$ and 0.038 dex in $`\mathrm{log}g`$. We adopt the same uncertainties in this analysis since both data sets are identical in terms of data acquisition, reduction, and S/N.
### 2.3 Results
#### 2.3.1 Adopted Atmospheric Parameters
The values of $`T_{\mathrm{eff}}`$ and $`\mathrm{log}g`$ for each of the 121 constant DA stars are listed in Table 1. We also include masses and absolute visual magnitudes derived from the evolutionary models of Wood (1995) with carbon-core compositions, helium layers of $`q(\mathrm{He})M_{\mathrm{He}}/M_{}=10^2`$, and thick hydrogen layers of $`q(\mathrm{H})=10^4`$. Several individual objects in Table 1 are worth discussing before looking at the global properties of the sample.
There are four known unresolved double degenerates included in our sample. The first three of those are G1-45 (WD 0101+048; Maxted et al., 2000), and LP 550-52 (WD 1022+050) and G21-15 (WD 1824+050; Maxted & Marsh, 1999). Liebert et al. (1991) have shown that in such cases the atmospheric parameters derived are in fact an average of the parameters of both components of the system. Similarly, Bergeron et al. (1990a) suggested on the basis of spectroscopic and energy distribution fits that GD 387 (WD 2003+437) is probably composed of a DA and a DC star. They derived $`T_{\mathrm{eff}}=14,340`$ K and $`\mathrm{log}g=7.50`$ for the DA component. Therefore, the atmospheric parameters reported here for these four systems are quite uncertain.
Three stars in Table 1 have composite spectra, and reprocessing of the EUV flux from the white dwarf primary in the chromosphere of the secondary contaminates the center of some, or all, of the Balmer lines. These are PG 0308+096 (Saffer et al., 1993), PG 1643+144 (Kidder, 1991) and Case 1 (WD 1213+528; Lanning, 1982). For PG 0308+096, the only contaminated line is H$`\beta `$. Therefore, we exclude that line from the fitting procedure and are able to get a satisfactory fit with atmospheric parameters identical to those reported in Table 2 of LBH. Similarly, in the case of PG 1643+144 we exclude both H$`\beta `$ and H$`\gamma `$. For Case 1 however, nearly all the spectral lines, and H$`\beta `$ in particular, are contaminated by the companion. As before, we exclude H$`\beta `$ but we also exclude 25 Å from either side of the line centers for H$`\gamma `$ through H$`ϵ`$, fitting only the line wings of the Balmer series. The effective temperature thus obtained, $`T_{\mathrm{eff}}=13,920`$ K agrees well enough with that determined by Sion et al. (1984) based on a fit of to the IUE spectrum, $`T_{\mathrm{eff}}=13,000\pm 500`$ K,
Finally, our sample also includes two stars known to be magnetic, GD 77 (WD 0637+477; Schmidt et al., 1992) and G128-72 (WD 2329+267; Moran et al., 1998). They both show the characteristic Zeeman splitting of the Balmer lines caused by their magnetic fields and thus fitting their spectra is problematic due to the additional spectral broadening. Therefore the atmospheric parameters reported here for these two objects remain uncertain. For instance we obtain for G128-72 a spectroscopic solution of $`T_{\mathrm{eff}}=11,520`$ K and $`\mathrm{log}g=9.09`$, while a fit to the $`BVRIJHK`$ photometric energy distribution combined with a trigonometric parallax measurement yields $`T_{\mathrm{eff}}=9400`$ K and $`\mathrm{log}g=8.02`$ according to Bergeron et al. (2001).
These uncertain atmospheric parameter measurements are indicated by colons in Table 1 and we must pay particular attention to the corresponding objects when discussing the ZZ Ceti instability strip below.
#### 2.3.2 G226-29
Before discussing the results of our analysis any further, we want to consider the case of G226-29. Being the hottest ZZ Ceti star analyzed by Bergeron et al. (2004) with $`T_{\mathrm{eff}}=12,460`$ K and $`\mathrm{log}g=8.28`$, G226-29 represents an important object for determining the slope of the blue edge of the ZZ Ceti instability strip (see Fig. 4 of Bergeron et al., 2004). These atmospheric parameter determinations are based on the same spectrum than the one used by B95 in their analysis of the atmospheric convective efficiency in DA white dwarfs. However, B95 also discuss a second spectroscopic observation of G226-29 with derived atmospheric parameters that agree within the uncertainties with the values given above. To be more specific, the atmospheric parameters derived from this second observation are $`T_{\mathrm{eff}}=12,260`$ K and $`\mathrm{log}g=8.32`$, consistent with the previous estimates within the uncertainties quoted in the previous section. What is more interesting perhaps is that this new temperature estimate is now in perfect agreement with the UV temperature obtained from the IUE spectrum, $`T_{\mathrm{eff}}=12,270`$ K (see Fig. 12 of B95). Given this improved internal consistency, we adopt from now on these new atmospheric parameters for G226-29. These are reported in Table 2 together with the values for the mass and absolute magnitude.
#### 2.3.3 New ZZ Ceti Stars
To complete the picture, in addition to the nonvariable stars given in Table 1, we need to include all ZZ Ceti stars for which we have spectroscopic observations. These include the 36 ZZ Ceti stars from Bergeron et al. (2004) as well as 3 new ZZ Ceti stars: PB 520 and GD 133 discovered by Silvotti et al. (2005) and Silvotti et al. (2005, in preparation), respectively, and G232-38 (WD 2148+539; $`V=16.4`$) discovered as part of our ongoing spectroscopic survey of the McCook & Sion catalog described in § 4. Our fits to the Balmer lines of these new variables are presented in Figure TOWARDS AN EMPIRICAL DETERMINATION OF THE ZZ CETI INSTABILITY STRIP; the atmospheric parameters for each object are reported in Table 2 together with the masses and absolute visual magnitudes. The values of $`T_{\mathrm{eff}}`$ and $`\mathrm{log}g`$ place PB 520 and G232-39 squarely within the limits of the ZZ Ceti instability strip (see Fig. TOWARDS AN EMPIRICAL DETERMINATION OF THE ZZ CETI INSTABILITY STRIP below) and we were more than confident that high speed photometric measurements would confirm their variability.
Silvotti et al. (2005) had already reported the detection of photometric variability in PB 520. G232-38, on the other hand, had never been observed before for photometric variability to our knowledge. Thus, we obtained high-speed photometric observations of G232-38 during an observing run in 2004 October at the 1.6 m telescope of the Observatoire du mont Mégantic equipped with LAPOUNE, the portable Montréal three-channel photometer. In all, we were able to obtain 3.9 h of data. Our sky-subtracted, extinction-corrected light curve of G232-38 is displayed in Figure TOWARDS AN EMPIRICAL DETERMINATION OF THE ZZ CETI INSTABILITY STRIP. G232-38 is clearly a ZZ Ceti star with multiperiodic luminosity variations. The resulting Fourier (amplitude) spectrum is displayed in Figure TOWARDS AN EMPIRICAL DETERMINATION OF THE ZZ CETI INSTABILITY STRIP. Three main-frequency components are easily discernible, with periods of 741.6, 984.0 and 1147.4 s. These relatively long periods and the rather large amplitude ($`10`$%) of the luminosity variations are consistent with the location of G232-38 somewhat closer to the red edge of the ZZ Ceti instability strip (see below).
After this paper was submitted, it has come to our attention that GD 133 (WD 1116+026) has been been identified as a short period ($`120`$ s), low-amplitude ($`<1`$%) ZZ Ceti star by Silvotti et al. (2005, in preparation) based on high-speed photometric observations obtained at the VLT with ULTRACAM. This object has long been thought to be photometrically constant according to numerous published sources (McGraw, 1977; Kepler et al., 1995; Giovannini, 1996; Silvotti et al., 1997). Back in March 2003, two of us (G.F. and P.B.) had even observed this star with the 61-inch telescope at the Mount Bigelow observatory, the light curve of which is displayed in Figure TOWARDS AN EMPIRICAL DETERMINATION OF THE ZZ CETI INSTABILITY STRIP. Although there is no obvious periodicity observed in the light curve, the corresponding Fourier (amplitude) spectrum shown in Figure TOWARDS AN EMPIRICAL DETERMINATION OF THE ZZ CETI INSTABILITY STRIP yields one significant peak above the 1$`\sigma `$ noise level with a period of 120.13 s, consistent with the observations of Silvotti et al. We have two independent optical spectra for GD 133, one from C. Moran (1999, private communication) that yields $`T_{\mathrm{eff}}=12,090`$ K and $`\mathrm{log}g=8.06`$, and our own data obtained in 2003 June that yields $`T_{\mathrm{eff}}=12,290`$ K and $`\mathrm{log}g=8.05`$. Although both sets of atmospheric parameters are consistent within the uncertainties, the former solution places GD 133 within the confines of our empirical instability strip and this is the solution we will adopt here. We report the atmospheric parameters for GD 133 in Table 2 along with our determination for the mass and absolute visual magnitude. We note that the location of GD 133 at the blue edge of the strip (see Fig. 6 below) is entirely consistent with the low amplitude and short pulsation period observed in Figures TOWARDS AN EMPIRICAL DETERMINATION OF THE ZZ CETI INSTABILITY STRIP and TOWARDS AN EMPIRICAL DETERMINATION OF THE ZZ CETI INSTABILITY STRIP .
#### 2.3.4 The Empirical ZZ Ceti Instability Strip
The locations of all 121 constant DA stars from Table 1 along with the 36 ZZ Ceti stars from Bergeron et al. (2004) and the 3 new ZZ Ceti stars discussed above are plotted in Figure TOWARDS AN EMPIRICAL DETERMINATION OF THE ZZ CETI INSTABILITY STRIP in a $`T_{\mathrm{eff}}\mathrm{log}g`$ diagram. Only 82 out the 121 nonvariables have atmospheric parameters that place them within the confines of Figure TOWARDS AN EMPIRICAL DETERMINATION OF THE ZZ CETI INSTABILITY STRIP. The bold open circles within the strip correspond, from left to right, to the new ZZ Ceti stars GD 133, PB 520, and G232-38.
Given this unbiased sample, we can clearly see that the ZZ Ceti stars define a trapezoidal region in the $`T_{\mathrm{eff}}`$-$`\mathrm{log}g`$ plane in which no nonvariable stars are found, within the measurement errors, in agreement with the conclusions of Bergeron et al. (2004) and references therein. And there is certainly no need here to go through any statistical analysis to conclude that the ZZ Ceti instability strip is indeed pure. We must also note that all nonvariable white dwarfs claimed to be close or even within the ZZ Ceti instability strip are in fact well outside the strip according to our analysis. These are GD 52 (WD 0348+339; Dolez et al., 1991; Silvotti et al., 1997); G8-8 (WD 0401+250; Silvotti et al., 1997; Kepler & Nelan, 1993); GD 31 (WD 0231$``$054), Rubin 70 (WD 0339+523), GD 202 (WD 1636+160) according to Dolez et al. (1991); PB 6089 (WD 0037$``$006) and G130-5 (WD 2341$`+`$322) according to Silvotti et al. (1997); BPM 20383 (WD 1053$``$550), BPM 2819 (WD 0255$``$705) according to Kepler & Nelan (1993, PG 1022$`+`$050 is a double degenerate); PG 1119+385, GD 515 (WD 1654+637), GD 236 (WD 2226+061) according to Kepler et al. (1995). There is also the case of GD 556 (WD 2311+552; Dolez et al., 1991; Kepler et al., 1995; Giovannini et al., 1998), which we find slightly hotter than the red edge of the strip; this object is discussed further in the next section.
One of the primary goals of our study is to improve the determination of the location of the blue and red edges of the empirical ZZ Ceti instability strip by using both variable and nonvariable DA white dwarfs. The results shown in Figure TOWARDS AN EMPIRICAL DETERMINATION OF THE ZZ CETI INSTABILITY STRIP first reveal that the location of the red edge is better constrained than the blue edge, in particular because of the three nonvariables (GD 556, GD 426, and EC 12043-1337) that lie very close to the red edge. In contrast, there are very few hot nonvariables near the blue edge. Note that the filled squares at the top of the figure are unresolved double degenerates and the atmospheric parameters obtained here are the average values of both components of the system. Hence these cannot be used to constrain the slope of the blue edge. In addition, our revised temperature for G226-29, which is 200 K cooler than our previous estimate, now removes the previous constraint we had on the slope of the blue edge. We show in Figure TOWARDS AN EMPIRICAL DETERMINATION OF THE ZZ CETI INSTABILITY STRIP the range of possibilities for the blue edge as defined by our spectroscopic analysis. It is clear that additional observations close to the blue edge are badly needed to constrain the slope better.
We point out, in this connection, that nonadiabatic pulsation theory does suggest that the slope of the blue edge in a $`T_{\mathrm{eff}}\mathrm{log}g`$ diagram such as the one shown in Figure TOWARDS AN EMPIRICAL DETERMINATION OF THE ZZ CETI INSTABILITY STRIP should be significantly smaller than that of the red edge, leading to an expected strip which is wider at higher surface gravities. The last word on the question of the theoretical ZZ Ceti instability strip has been presented by Fontaine et al. (2003). We show in Figure TOWARDS AN EMPIRICAL DETERMINATION OF THE ZZ CETI INSTABILITY STRIP an updated comparison with their theoretical results (solid lines). We find that the slope of the theoretical blue edge is compatible with the range of possibilities allowed by our empirical results. On the other hand, the slope of the theoretical red edge is not too different from our own determination however it is predicted to be somewhat hotter than the red edge inferred from observation. Our aim in the future is to focus on the $`empirical`$ boundaries with improved statistics, especially for the blue edge.
A global characteristic that is also noticeable in Figure TOWARDS AN EMPIRICAL DETERMINATION OF THE ZZ CETI INSTABILITY STRIP is the trend toward higher values of $`\mathrm{log}g`$ as $`T_{\mathrm{eff}}`$ decreases. This is now a familiar result observed in all spectroscopic surveys extending to low temperatures (B95, Koester et al. 2001, Kleinman et al. 2004, LBH, Gianninas et al. 2005). It has been proposed by Bergeron et al. (1990b) that these high inferred masses could be the result of small amounts of helium brought to the surface by the hydrogen convection zone, hence increasing the atmospheric pressure. When analyzed with pure hydrogen models, this increased pressure could be misinterpreted as resulting from of a high mass (see also Boudreault & Bergeron, 2005).
#### 2.3.5 GD 556
One constant star in Figure TOWARDS AN EMPIRICAL DETERMINATION OF THE ZZ CETI INSTABILITY STRIP, GD 556, has an effective temperature slightly hotter than the empirical red edge. If we refer to Table 1, there are four independent sources that concluded that GD 556 is not a variable DA white dwarf. However, we would like to recall that initially G30-20 had also been found to be constant by Dolez et al. (1991) and Bergeron & McGraw (1989, unpublished) but was later identified as a ZZ Ceti pulsator by Mukadam et al. (2002); GD 133 discussed above is also a good example. It is worth mentioning that GD 556 presents certain challenges as far as photometric observations are concerned. Firstly, it is a rather dim star with $`V16.2`$ (McCook & Sion, 1999). Secondly, its position near the red edge implies that if it is indeed a pulsator, it should show long period pulsations that can be difficult to detect if one observes the star while two pulsational modes are interfering destructively. Considering all these facts, we believe that GD 556 is definitely worth re-observing under favorable conditions, both photometrically and spectroscopically. Nonetheless, if GD 556 truly is photometrically constant, then considering our error bars, the fact that it lies within the strip, albeit very close to the red edge, changes nothing in our conclusions relative to the the purity of the ZZ Ceti instability strip.
## 3 RESULTS FROM THE SLOAN DIGITAL SKY SURVEY
Since the discovery of the first pulsating DA white dwarf by Landolt (1968), HL Tau 76, and up to the spectroscopic study of Bergeron et al. (2004), a total of 36 ZZ Ceti stars were known (see Table 1 of Bergeron et al., 2004), a quarter of which had been discovered using the spectroscopic technique. In a single effort, Mukadam et al. (2004a) reported the discovery of 35 new ZZ Ceti pulsators, hence nearly doubling the number of known variables in this class. Thirty three of these have been discovered in the white dwarf SDSS sample, mostly from the first data release (Kleinman et al., 2004), and two more from the Hamburg Quasar Survey. Very recently, Mullally et al. (2005) reported the discovery of eleven more ZZ Ceti stars from SDSS as well as several nonvariable stars. However, the stars from Mullally et al. (2005) are not included in the analysis and discussion that follow.
ZZ Ceti candidates from the SDSS were selected for follow-up high-speed photometry on the basis of various techniques including $`ugriz`$ photometry, equivalent width measurements, and the spectroscopic technique using SDSS spectra and Koester’s model atmospheres. By far, the spectroscopic technique led to a significantly higher success rate of discovery than other techniques (90% by confining the candidates between $`T_{\mathrm{eff}}=11,000`$ K and 12,000 K). The 33 new SDSS pulsators are listed in Table 1 of Mukadam et al. (2004a), while nonvariables are given in their Tables 2 and 3 for different detection thresholds. An examination of these tables reveal that all objects are relatively faint ($`g17`$) due to the intrinsic characteristics of the Sloan survey, which is aimed at identifying distant galaxies and quasars. Stellar objects on a given plate with an assigned fiber had to be faint in order not to saturate the detector.
Even though effective temperatures and surface gravities obtained from spectroscopic fits were provided in their paper, Mukadam et al. (2004a) did not discuss the implications of their new discoveries on the empirical determination of the ZZ Ceti instability strip. That discussion was deferred to a second paper by Mukadam et al. (2004b) who analyzed in more detail the spectroscopic results from their first paper, with a particular emphasis on the empirical ZZ Ceti instability strip as inferred from the location of variables and nonvariables in the $`T_{\mathrm{eff}}\mathrm{log}g`$ plane. In particular, the authors of that study question one more time the purity of the ZZ Ceti instability strip. The results of their analysis are contrasted with our results in Figure TOWARDS AN EMPIRICAL DETERMINATION OF THE ZZ CETI INSTABILITY STRIP. We should mention that both analyses rely on different sets of model atmospheres (ours versus D. Koester’s models) and there could be systematic offsets. But the most striking feature of the Mukadam et al. results is the large number of nonvariable white dwarfs within their empirical instability strip.
Through a painstaking statistical analysis of their results, Mukadam et al. (2004b) conclude that 18 nonvariables fall within the ZZ Ceti instability strip. Given that 33 new pulsators have been discovered from the same sample, the results suggest that the ZZ Ceti instability strip is only $`50`$% pure, at best. The authors have even estimated the probability that the instability strip is pure is only 0.004 %! This result is of course in sharp contrast with our conclusions based on a comparable number of white dwarfs, and considerably brighter than those discovered in the SDSS. If indeed the instability strip is contaminated by a significant fraction of nonvariables, as implied by Mukadam et al., then the global properties of DA stars inferred from asteroseismological studies of ZZ Ceti stars could not be generalized to the entire population of hydrogen-atmosphere white dwarfs as the ZZ Ceti pulsators would no longer represent a phase through which all DA stars must evolve. Another important implication of this challenging result is that the pulsation instability of a white dwarf would no longer depend solely on its effective temperature and stellar mass, but would require an additional, yet unidentified, physical parameter to discriminate variables and nonvariables within the instability strip.
How can our results be reconciled with those of Mukadam et al.? The authors claim that since the discovery of white dwarf variables in 1968, their study represents the first analysis of a homogeneous set of spectra acquired using the same instrument on the same telescope, and with consistent data reductions. There is even an implicit suggestion that this homogeneity could account for the fundamental difference between their analysis and that of Bergeron et al. (2004). However, this point of view completely ignores the incentive behind the earlier study of B95 whose specific goal was to provide an analysis of a homogeneous set of spectroscopic observations of the 18 ZZ Ceti stars known at that time, observable from the northern hemisphere. As discussed in § 2.1 and 2.2 of B95, the first spectroscopic analysis of a sizeable sample of ZZ Ceti stars by Daou et al. (1990) relied on spectra acquired as part of a backup project by various observers, and thus with different telescopes, spectrographs, detectors, and reduction procedures. As such, the spectroscopic sample of Daou et al. was somewhat inhomogeneous. To overcome precisely this problem, it was deemed necessary for B95 to reacquire optical spectra for the ZZ Ceti stars using the same instrument setup and reduction techniques. Hence high S/N spectroscopic observations for the 18 ZZ Ceti stars were acquired using the 2.3-m telescope at Steward Observatory, equipped with the Boller & Chivens spectrograph and a Texas Instrument CCD detector; spectra of four additional ZZ Ceti stars from the southern hemisphere have also been analyzed by Bergeron et al., but even though the spectra were of comparable quality to those obtained at the Steward Observatory, these four stars were treated separately throughout their analysis to preserve the homogeneity of the spectroscopic sample.
Note that the same instrument setup has been used ever since in many of our studies, and in particular in the recent extensive spectroscopic analysis of LBH who reported effective temperatures and surface gravities for nearly 350 DA stars drawn from the Palomar Green (PG) Survey. Even though the completeness of the PG survey remains questionable, the sample analyzed by Liebert et al. represents one of the largest statistically significant samples of DA stars analyzed to date. Yet, only one ZZ Ceti candidate (PG 1349$`+`$552) was found within the empirical instability strip together with 9 previously known ZZ Ceti stars. High-speed photometric observations by Bergeron et al. (2004) confirmed that PG 1349$`+`$552 was indeed a new ZZ Ceti pulsator. Hence the conclusions of LBH are consistent with those of Bergeron et al. (2004), with the results presented in this paper, and with the results of our ongoing survey of the McCook & Sion catalog discussed in § 4, that the empirical instability strip contains no nonvariable stars. Hence, arguments based on the homogeneity of the spectroscopic analyses are unlikely to be able to explain the discrepancy between our conclusions and the contrasting results of Mukadam et al. (2004b). If anything, a spectroscopic analysis of an inhomogeneous data set should lead to a contamination of the instability strip with nonvariables, not the other way around!
Mukadam et al. (2004b) also suggested that their analysis effectively samples a different population of stars, more distant by a factor of 10 than that of the Bergeron et al. (2004) sample. Actually, taking a median value of $`g18.5`$ (see Fig. TOWARDS AN EMPIRICAL DETERMINATION OF THE ZZ CETI INSTABILITY STRIP below) and an absolute magnitude of $`M_g=11.64`$ obtained from a model atmosphere at $`T_{\mathrm{eff}}=12,000`$ K and $`\mathrm{log}g=8`$, we derive a distance of only 230 pc, still relatively close by. There is really no astrophysical reason to expect white dwarfs at that distance to behave differently from those at shorter distances. Other explanations must thus be sought.
A close examination of the 18 nonvariables claimed to be within the instability strip by Mukadam et al. (2004b, see their Table 1) reveals that all objects are among the faintest in their SDSS sample, as can be seen from Figure TOWARDS AN EMPIRICAL DETERMINATION OF THE ZZ CETI INSTABILITY STRIP where the distribution of SDSS white dwarfs taken from Tables 1 to 3 of Mukadam et al. (2004a) is shown as a function of the $`g`$ magnitude in the $`ugriz`$ photometric system. As discussed by B95, the S/N of the spectroscopic observations is one of the key aspects of the spectroscopic technique for determining precise atmospheric parameters, the other important one being the flux calibration. Since the exposure time of a given SDSS spectrum is set by that of the entire plate, the corresponding S/N must necessarily be a function of the magnitude of the star. To verify this assertion, we have measured the S/N values<sup>1</sup><sup>1</sup>1Here the S/N is measured in the continuum between 4450 and 4750 Å. of all SDSS spectra taken from Mukadam et al. (2004a), and plotted these values against the corresponding $`g`$ magnitude. This is shown in Figure TOWARDS AN EMPIRICAL DETERMINATION OF THE ZZ CETI INSTABILITY STRIP. As expected, fainter stars have lower S/N spectra, and only objects brighter than $`g17`$ have S/N above 40. This is not the case with standard slit spectroscopy, however, where the exposure time can be adjusted on a star-to-star basis. In B95 for instance, the exposure times were set to achieve an imposed lower limit of $`\mathrm{S}/\mathrm{N}80`$, although most spectra had $`\mathrm{S}/\mathrm{N}100`$ since the exposure times were also set long enough to cover several pulsation cycles (for an average of $`4.8`$ cycles) in order to obtain meaningful time-averaged spectra. We mention that this last criterion is not necessarily met in the SDSS spectroscopic data.
We now turn to a more detailed comparison of S/N between our spectroscopic sample and that of Mukadam et al. (2004a). We show in the top panel of Figure TOWARDS AN EMPIRICAL DETERMINATION OF THE ZZ CETI INSTABILITY STRIP the distribution of S/N values for our spectroscopic sample, including the photometrically constant stars from Table 1, the 36 ZZ Ceti stars from Bergeron et al. (2004), and the 3 new pulsators from Table 2. Out of 39 ZZ Ceti stars, 12 have spectra with an admittedly lower S/N value than the imposed lower limit of $`80`$ set by B95 in their analysis. These spectra correspond to data provided to us by C. Moran (1999, private communication) in the course of his search for double degenerate binaries, or to spectroscopic observations obtained prior to the discovery of the photometric variability of the object; this includes the two ZZ Ceti stars PB 520 and G232-38 analyzed in this paper. Still, only two ZZ Ceti stars have spectra with $`\mathrm{S}/\mathrm{N}<50`$, and none below 40. The spectra for our photometrically constant sample also have fairly high S/N values, almost all above 50. In contrast, the S/N of the SDSS spectra<sup>2</sup><sup>2</sup>2Only a fraction of the SDSS spectra could be recovered from the SDSS Web site. shown in the bottom panel of Figure TOWARDS AN EMPIRICAL DETERMINATION OF THE ZZ CETI INSTABILITY STRIP have considerably lower values, with most spectra having $`\mathrm{S}/\mathrm{N}<60`$. Even worse, the subsample of nonvariables claimed to lie within the instability strip (hatched histogram) has even lower S/N values, with most objects having $`\mathrm{S}/\mathrm{N}<40`$. Hence, it is perhaps not too surprising that the results of Mukadam et al. (2004b) regarding the purity of the ZZ Ceti instability strip, which are based on an analysis of low S/N spectra, differ so much from our own conclusions based on much higher quality spectroscopic observations.
We illustrate in Figure TOWARDS AN EMPIRICAL DETERMINATION OF THE ZZ CETI INSTABILITY STRIP a typical spectrum from our own sample of ZZ Ceti stars, GD 66 with $`\mathrm{S}/\mathrm{N}=80`$, and one from the SDSS sample, SDSS J084746.81$`+`$451006.3 with $`\mathrm{S}/\mathrm{N}=20`$. The S/N value of the latter is actually more typical of the sample of nonvariables found within the strip (bottom panel of Fig. TOWARDS AN EMPIRICAL DETERMINATION OF THE ZZ CETI INSTABILITY STRIP). It is clear that the spectroscopic solution will necessarily depend on the quality of these spectra. To quantify this assertion, we performed a Monte Carlo simulation by taking a series of model spectra at $`T_{\mathrm{eff}}=12,000`$ K and $`\mathrm{log}g=8.0`$, by adding random noise to achieve a given signal-to-noise ratio, and by fitting these spectra with our standard fitting procedure. The resulting $`T_{\mathrm{eff}}`$ and $`\mathrm{log}g`$ values are then used to compute the standard deviations $`\sigma _{T_{\mathrm{eff}}}`$ and $`\sigma _{\mathrm{log}g}`$ for this assumed S/N value. Values of S/N from 10 to 200 were explored, thus encompassing the entire range exhibited by the spectra analyzed in this paper and by Mukadam et al. (2004b). The results of this exercise are displayed in Figure TOWARDS AN EMPIRICAL DETERMINATION OF THE ZZ CETI INSTABILITY STRIP. It is clear that stars with low S/N spectra will yield atmospheric parameters with larger internal uncertainties than those derived from higher quality observations. In particular, if we again take S/N = 20 as indicative of the SDSS stars, we see from Figure TOWARDS AN EMPIRICAL DETERMINATION OF THE ZZ CETI INSTABILITY STRIP that such spectra would yield effective temperatures that are uncertain by $`500`$ K. If we consider that the width of the empirical instability strip is $`1000`$ K, it is easy to understand how lower quality spectra could easily place non-variable stars within the strip and vice-versa. Furthermore, stars with S/N $`80`$, typical of our photometric sample, exhibit uncertainties of roughly 150 K, which is entirely consistent with the uncertainties quoted in §2.2. Thus despite the homogeneous characteristics of the SDSS spectra in terms of instrument, telescope, and data reductions, their typical S/N is most likely too low to allow a precise measurement of the atmospheric parameters for these stars, or to determine accurately the location of the empirical ZZ Ceti instability strip, or to assess the purity of the strip for that matter.
Finally, we examine in Figure TOWARDS AN EMPIRICAL DETERMINATION OF THE ZZ CETI INSTABILITY STRIP the location of the SDSS white dwarfs in the $`T_{\mathrm{eff}}\mathrm{log}g`$ plane. In each panel we consider only the objects with spectra above a certain threshold in S/N (the bottom panel includes all objects). Also reproduced is the empirical ZZ Ceti instability strip determined by Bergeron et al. (2004). The top panel with $`\mathrm{S}/\mathrm{N}>70`$ corresponds to a threshold that would include $`80`$% of all ZZ Ceti stars from the sample of Bergeron et al. (2004, top panel of Fig. TOWARDS AN EMPIRICAL DETERMINATION OF THE ZZ CETI INSTABILITY STRIP). By comparison, only 2 objects from the SDSS sample meet this criterion. Hence if we restrict the analysis to the best spectra of both samples, the results are consistent: all variables are found within the empirical strip and all nonvariables lie outside. For $`\mathrm{S}/\mathrm{N}>40`$ (middle panel), 9 objects from the SDSS sample are found in the temperature range shown in Figure TOWARDS AN EMPIRICAL DETERMINATION OF THE ZZ CETI INSTABILITY STRIP. In this case, however, one variable star (WD 1711+6541 at $`T_{\mathrm{eff}}=11,310`$ K and $`\mathrm{log}g=8.64`$) falls slightly below the empirical red edge of the strip, while one nonvariable star (WD 1338$``$0023 at $`T_{\mathrm{eff}}=11,650`$ K and $`\mathrm{log}g=8.08`$) sits comfortably near the middle of the strip. Since these two objects are relatively bright ($`g=16.89`$ and 17.09, respectively), we managed to secure our own spectroscopic observations of these stars using the Steward Observatory 2.3 m telescope during an observing run in 2004 May. Reassuringly enough, our independent analysis of these two objects in terms of both data and models — $`T_{\mathrm{eff}}=11,490`$ K and $`\mathrm{log}g=8.56`$ for WD 1711+6541 and $`T_{\mathrm{eff}}=11,980`$ K and $`\mathrm{log}g=7.94`$ for WD 1338$``$0023 — places them where they are expected, that is, inside and outside the instability strip, respectively (see Fig. TOWARDS AN EMPIRICAL DETERMINATION OF THE ZZ CETI INSTABILITY STRIP).
For completeness, we show at the bottom of Figure TOWARDS AN EMPIRICAL DETERMINATION OF THE ZZ CETI INSTABILITY STRIP all the objects from the SDSS sample ($`\mathrm{S}/\mathrm{N}>0`$). Once again, we can see that the bulk of this sample is characterized with S/N values below 40, the threshold value used in the middle panel, and that the conclusion about the purity of the ZZ Ceti instability strip rests heavily on the quality of the spectroscopic observations.
## 4 ONGOING SPECTROSCOPIC SURVEY
In order to increase the number of stars in our spectroscopic (and eventually photometric) sample, we have undertaken a broader spectroscopic survey of DA stars drawn from the Catalog of Spectroscopically Identified White Dwarfs of McCook & Sion (1999). We have defined our sample using the following criteria: (1) a temperature index lying between 3 and 7, (2) apparent visual magnitudes of V $`<`$ 17 and (3) declination greater than $``$30 degrees. High S/N optical spectroscopic observations are currently being secured for each star that meets these criteria. This survey was initiated with several goals in mind. First and foremost, we wish to obtain measurements of $`T_{\mathrm{eff}}`$ and $`\mathrm{log}g`$ for each object. Secondly, we want to confirm the spectroscopic classification of stars from the catalog (we have already identified 29 stars misclassified as DA stars that are clearly lower gravity objects). A final goal of our survey is to identify new ZZ Ceti candidates (G232-38 has been discovered in this survey). This is the reason for restricting ourselves to stars with the aforementioned range of temperature indices. Some preliminary results of this analysis have already been presented in Gianninas et al. (2005). Among these is the discovery of a unique DAZ white dwarf, GD 362 (Gianninas et al., 2004).
The combined results of our ongoing spectroscopic survey and of the photometric sample analyzed in § 2 are displayed in Figure TOWARDS AN EMPIRICAL DETERMINATION OF THE ZZ CETI INSTABILITY STRIP as triangles and circles, respectively. The three low-gravity objects in the vicinity of the instability strip are known double degenerate systems. We have already discussed two of these (see § 2.3.4 and Fig. TOWARDS AN EMPIRICAL DETERMINATION OF THE ZZ CETI INSTABILITY STRIP), the third is GD 429 (Maxted et al., 2000) which has yet to be observed for photometric variability. We clearly see that many objects from our survey lie very close to both the red and blue edges of the instability strip. These stars are important as we attempt to determine better the exact boundaries of the instability strip. Therefore, we plan on securing high speed photometric observations for these objects in order to confirm their photometric status. These results will be reported in due time.
## 5 CONCLUSION
We have gathered optical spectra for 121 photometrically constant DA white dwarfs for which we derived values of $`T_{\mathrm{eff}}`$ and $`\mathrm{log}g`$. Using these nonvariable white dwarfs together with a sample of 39 relatively bright ZZ Ceti stars, we wished to obtain a better understanding of the location and shape of the red and blue edges of the ZZ Ceti instability strip. In so doing, we have succeeded in better populating the $`T_{\mathrm{eff}}`$-$`\mathrm{log}g`$ plane in the vicinity of the ZZ Ceti instability strip. We find that the location and slope of the red edge is quite well constrained whereas our newly adopted atmospheric parameters for G226-29 allow for a much broader range of slopes for the blue edge which would accommodate our current photometric sample. Furthermore, we find no nonvariable white dwarfs within the ZZ Ceti instability. This supports our belief that ZZ Ceti stars represent an evolutionary stage by which all DA white dwarfs must pass.
The optical spectra that we analyzed were gathered as part of a more extensive survey of DA white dwarfs from the catalog of McCook & Sion (1999). This survey has several goals, among them, the identification of candidate ZZ Ceti stars. Thus far, two of these, PB 520 and G232-38, have been identified as ZZ Ceti pulsators by Silvotti et al. (2005) and us, respectively. The spectroscopic technique pioneered by B95 has proven once again to be an invaluable tool as far as identifying new candidate ZZ Ceti stars. Indeed, it has maintained its 100% success rate in predicting variability in DA white dwarfs. With the inclusion of PB 520, G232-38, and GD 133 the number of ZZ Ceti stars (excluding the SDSS stars) swells to 39 of which 11 have been successfully identified using this method. Even among the ZZ Ceti stars discovered through SDSS, the most fruitful method for identifying candidates was the spectroscopic method (Mukadam et al., 2004a). In the future, in order to define better the blue edge, and to study further the instability strip as a whole, we plan on securing high speed photometric observations for all the DA white dwarfs that are within or near the boundaries of the empirical strip and that have never been observed for variability.
We have also been able to show the importance of using high-quality data (i.e., high S/N) when performing analyses such as these through an in-depth examination of the data used by Mukadam et al. (2004b) in their study. It is clear that their controversial results, which place a large number of nonvariable stars within the instability strip, can be traced back to spectra of lesser quality that greatly affect the result of the spectroscopic fit. However, one cannot discount the fact that Mukadam et al. (2004a) have nearly doubled the number of known ZZ Ceti stars as well as adding a large number of nonvariable DA white dwarfs to the mix. The study of the ZZ Ceti instability strip can only benefit from the inclusion of all recently identified variable and nonvariable DA white dwarfs within our sample. We are therefore exploring the possibility of obtaining high S/N optical spectra for all the DA white dwarfs from Mukadam et al. (2004a), Silvotti et al. (2005), and Mullally et al. (2005) in the near future.
We would like to thank the director and staff of Steward Observatory for the use of their facilities. We would also like to thank the director and staff of the Observatoire du mont Mégantic for the use of their facilities and for supporting LAPOUNE as a visitor instrument. We also acknowledge the contribution of F. Provencher in the analysis of the SDSS data. This work was supported in part by the NSERC Canada and by the Fonds Québécois de la recherche sur la nature et les technologies (Québec). GF acknowledges the contribution of the Canada Research Chair Program. |
warning/0506/math0506215.html | ar5iv | text | # Scaled Enflo type is equivalent to Rademacher type
## 1 Introduction
Recall that a Banach space $`X`$ is said to have Rademacher type $`p>0`$ (see ) if there exists a constant $`T<\mathrm{}`$ such that for every $`x_1,\mathrm{},x_nX`$,
$`𝔼_\epsilon {\displaystyle \underset{j=1}{\overset{n}{}}}\epsilon _jx_j_X^pT^p{\displaystyle \underset{j=1}{\overset{n}{}}}x_j_X^p,`$ (1)
where here, and in what follows, $`𝔼_\epsilon `$ denotes the expectation with respect to uniformly chosen $`\epsilon =(\epsilon _1,\mathrm{},\epsilon _n)\{1,1\}^n`$. The infimum over all constants $`T`$ for which (1) holds is denoted $`T_p(X)`$.
Motivated by the search for concrete versions of Ribe’s theorem for various fundamental local properties of Banach spaces (see the discussion in ), several researchers proposed non-linear notions of type, which make sense in the setting arbitrary metric spaces (see ). In particular, following Enflo we say that a metric space $`(,d_{})`$ has Enflo type $`p`$ if there exists a constant $`K`$ such that for every $`n`$ and every $`f:\{1,1\}^n`$,
$`𝔼_\epsilon d_{}(f(\epsilon ),f(\epsilon ))^pT^p{\displaystyle \underset{j=1}{\overset{n}{}}}𝔼_\epsilon d_{}(f(\epsilon _1,\mathrm{},\epsilon _{j1},\epsilon _j,\epsilon _{j+1},\mathrm{},\epsilon _n),f(\epsilon _1,\mathrm{},\epsilon _{j1},\epsilon _j,\epsilon _{j+1},\mathrm{},\epsilon _n))^p.`$ (2)
For Banach spaces (1) follows from (2) by considering the function $`\epsilon _{j=1}^n\epsilon _jx_j`$. The question whether in the category of Banach spaces Rademacher type $`p`$ implies Enflo type $`p`$ was posed by Enflo in , and in full generality remains open. In Pisier showed that if a Banach space has Rademacher $`p`$ then it has Enflo type $`p^{}`$ for every $`p^{}<p`$ (see also the work of Bourgain, Milman and Wolfson for a similar result which holds for a another notion of non-linear type). In it was shown that for UMD Banach spaces (see ) Rademacher type $`p`$ is equivalent to Enflo type $`p`$.
Motivated by our recent work on metric cotype , we introduce below the notion of scaled Enflo type of a metric space (which is, in a sense, “opposite” to the notion of metric cotype defined in ), and show that for Banach spaces, scaled Enflo type $`p`$ is equivalent to Rademacher type $`p`$. This settles the long standing problem of finding a purely metric formulation of the notion of type (though Enflo’s problem described above remains open). Modulo some of the results of , the proof of our main theorem is very simple.
###### Definition 1.1 (Scaled Enflo type).
Let $`(,d_{})`$ be a metric space and $`p>0`$. We say that $``$ has scaled Enflo type $`p`$ with constant $`\tau `$ if for every integer $`n`$ there exists an even integer $`m`$ such that for every $`f:_m^n`$,
$`𝔼_\epsilon {\displaystyle _{_m^n}}d_{}(f\left(x+{\displaystyle \frac{m}{2}}\epsilon \right),f(x))^p𝑑\mu (x)\tau ^pm^p{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle _{_m^n}}d_{}(f(x+e_j),f(x))^p𝑑\mu (x),`$ (3)
where $`\mu `$ is the uniform probability measure on $`_m^n`$, and $`\{e_j\}_{j=1}^n`$ is the standard basis of $`^n`$. The infimum over all constants $`\tau `$ for which (3) holds is denoted $`\tau _p()`$.
###### Theorem 1.2.
Let $`X`$ be a Banach space and $`p[1,2]`$. Then $`X`$ has Rademacher type $`p`$ if and only if $`X`$ has scaled Enflo type $`p`$. More precisely,
$$\frac{1}{2\pi }T_p(X)\tau _p(X)5T_p(X).$$
## 2 Proof of Theorem 1.2
We start by showing that scaled Enflo type $`p`$ implies Rademacher type $`p`$.
###### Lemma 2.1.
Let $`X`$ be a Banach space and $`p[1,2]`$. Then $`T_p(X)2\pi \tau _p(X)`$.
###### Proof.
Let $`X`$ be a Banach space and assume that $`\tau _p(X)<\mathrm{}`$ for some $`p[1,2]`$. Fix $`\tau >\tau _p(X)`$, $`v_1,\mathrm{},v_nX`$, and let $`m`$ be an even integer. Define $`f:_m^nX`$ by $`f(x_1,\mathrm{},x_n)=_{j=1}^ne^{\frac{2\pi ix_j}{m}}v_j`$. Then
$`{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle _{_m^n}}f\left(x+e_j\right)f(x)_X^p𝑑\mu (x)=\left|e^{\frac{2\pi i}{m}}1\right|^p{\displaystyle \underset{j=1}{\overset{n}{}}}v_j_X^p\left({\displaystyle \frac{2\pi }{m}}\right)^p{\displaystyle \underset{j=1}{\overset{n}{}}}v_j_X^p,`$ (4)
and
$`𝔼_\epsilon {\displaystyle _{_m^n}}f\left(x+{\displaystyle \frac{m}{2}}\epsilon \right)f(x)_X^p𝑑\mu (x)=2^p{\displaystyle _{_m^n}}{\displaystyle \underset{j=1}{\overset{n}{}}}e^{\frac{2\pi ix_j}{m}}v_j_X^p𝑑\mu (x).`$ (5)
We recall the contraction principle (see ), which states that for every $`a_1,\mathrm{},a_n`$,
$$𝔼_\epsilon \underset{j=1}{\overset{n}{}}\epsilon _ja_jv_j_X^p(\underset{1jn}{\mathrm{max}}|a_j|)^p𝔼_\epsilon \underset{j=1}{\overset{n}{}}\epsilon _jv_j_X^p.$$
Thus,
$`{\displaystyle _{_m^n}}{\displaystyle \underset{j=1}{\overset{n}{}}}e^{\frac{2\pi ix_j}{m}}v_j_X^pd\mu (x)={\displaystyle _{_m^n}}𝔼_\epsilon {\displaystyle \underset{j=1}{\overset{n}{}}}e^{\frac{2\pi i}{m}\left(x_j+\frac{m(1\epsilon _j)}{4}\right)}v_j_X^pd\mu (x)={\displaystyle _{_m^n}}𝔼_\epsilon {\displaystyle \underset{j=1}{\overset{n}{}}}\epsilon _je^{\frac{2\pi ix_j}{m}}v_j_X^pd\mu (x){\displaystyle \frac{1}{2^p}}𝔼_\epsilon {\displaystyle \underset{j=1}{\overset{n}{}}}\epsilon _jv_j_X^p.`$ (6)
Combining (4), (5) and (6) yields the required result. ∎
Let $`X`$ be a Banach space with type $`p`$, $`m`$ an integer divisible by $`4`$, and $`k`$ an odd integer. Fix $`f:_m^nX`$ and $`\epsilon \{1,1\}^n`$. Define $`𝒜^{(k)}f:_m^nX`$ by
$$𝒜^{(k)}f(x)=\frac{1}{k^n}\underset{z(k,k)^n(2)^n}{}f(x+z).$$
###### Lemma 2.2.
For $`p1`$ and every $`f:_m^nX`$
$$_{_m^n}𝒜^{(k)}f(x)f(x)_X^p𝑑\mu (x)(k1)^pn^{p1}\underset{j=1}{\overset{n}{}}_{_m^n}f(x+e_j)f(x)_X^p𝑑\mu (x).$$
###### Proof.
For every $`t`$ let $`s(t)`$ be the sign of $`t`$ (with convention that $`s(0)=0`$). For every $`z_m^n`$,
$`f(x+z)f(x)_X^pz_1^{p1}{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle \underset{\mathrm{}=1}{\overset{|z_j|}{}}}f\left(x+{\displaystyle \underset{t=1}{\overset{j1}{}}}z_te_t+\mathrm{}s(z_j)e_j\right)f\left(x+{\displaystyle \underset{t=1}{\overset{j1}{}}}z_te_t+(\mathrm{}1)s(z_j)e_j\right)_X^p.`$
Observe that since $`k`$ is odd, $`|(k,k)^n(2)^n|=k^n`$. Thus
$`{\displaystyle _{_m^n}}𝒜^{(k)}f(x)f(x)_X^p𝑑\mu (x){\displaystyle \frac{1}{k^n}}{\displaystyle \underset{z(k,k)^n(2)^n}{}}{\displaystyle _{_m^n}}f(x+z)f(x)_X^p𝑑\mu (x)`$
$``$ $`{\displaystyle \frac{1}{k^n}}{\displaystyle \underset{z(k,k)^n(2)^n}{}}{\displaystyle _{_m^n}}z_1^{p1}{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle \underset{\mathrm{}=1}{\overset{|z_j|}{}}}f\left(x+{\displaystyle \underset{t=1}{\overset{j1}{}}}z_te_t+\mathrm{}s(z_j)e_j\right)f\left(x+{\displaystyle \underset{t=1}{\overset{j1}{}}}z_te_t+(\mathrm{}1)s(z_j)e_j\right)_X^pd\mu (x)`$
$``$ $`{\displaystyle \frac{1}{k^n}}{\displaystyle \underset{z(k,k)^n(2)^n}{}}{\displaystyle \underset{j=1}{\overset{n}{}}}z_1^{p1}|z_j|{\displaystyle _{_m^n}}f(y+s(z_j)e_j)f(y)_X^p𝑑\mu (x)`$
$``$ $`(k1)^pn^{p1}{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle _{_m^n}}f(x+e_j)f(x)_X^p𝑑\mu (x).`$
###### Proof of theorem 1.2.
Fix an odd integer $`k`$, with $`k<\frac{m}{2}`$. As in , given $`j\{1,\mathrm{},n\}`$ we define $`S(j,k)_m^n`$ by
$$S(j,k)\{x[k,k]^n_m^n:y_j0mod2\mathrm{and}\mathrm{}j,y_{\mathrm{}}1mod2\}.$$
For $`f:_m^nX`$ we define
$`_j^{(k)}f(x)=\left(f{\displaystyle \frac{\mathrm{𝟏}_{S(j,k)}}{\mu (S(j,k))}}\right)(x)={\displaystyle \frac{1}{\mu (S(j,k))}}{\displaystyle _{S(j,k)}}f(x+y)𝑑\mu (y).`$ (7)
In (see equation (39) there) it is shown that for every $`x_m^n`$ and $`\epsilon \{1,1\}^n`$,
$$\left(\frac{k}{k+1}\right)^{n1}\left(𝒜^{(k)}f(x+\epsilon )𝒜^{(k)}f(x\epsilon )\right)=\underset{j=1}{\overset{n}{}}\epsilon _j\left[_j^{(k)}f(x+e_j)_j^{(k)}f(xe_j)\right]+U(x,\epsilon )+V(x,\epsilon ),$$
where, by inequalities (41) and (42) in , for every $`\epsilon \{1,1\}^n`$,
$$\mathrm{max}\{_{_m^n}U(x,\epsilon )_X^p𝑑\mu (x),_{_m^n}V(x,\epsilon )_X^p𝑑\mu (x)\}\frac{8^pn^{2p1}}{k^p}\underset{j=1}{\overset{n}{}}_{_m^n}f(x+e_j)f(x)_X^p.$$
Thus, for every $`T>T_p(X)`$,
$`\left({\displaystyle \frac{k}{k+1}}\right)^{p(n1)}𝔼_\epsilon {\displaystyle _{_m^n}}𝒜^{(k)}f(x+\epsilon )𝒜^{(k)}f(x\epsilon )_X^p𝑑\mu (x)`$ (8)
$``$ $`3^{p1}𝔼_\epsilon {\displaystyle _{_m^n}}{\displaystyle \underset{j=1}{\overset{n}{}}}\epsilon _j\left[_j^{(k)}f(x+e_j)_j^{(k)}f(xe_j)\right]_X^p𝑑\mu (x)+{\displaystyle \frac{24^pn^{2p1}}{k^p}}{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle _{_m^n}}f(x+e_j)f(x)_X^p𝑑\mu (x)`$
$``$ $`3^{p1}T^p{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle _{_m^n}}_j^{(k)}f(x+e_j)_j^{(k)}f(xe_j)_X^p𝑑\mu (x)+{\displaystyle \frac{24^pn^{2p1}}{k^p}}{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle _{_m^n}}f(x+e_j)f(x)_X^p𝑑\mu (x)`$
$``$ $`3^{p1}T^p{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle _{_m^n}}f(x+e_j)f(xe_j)_X^p𝑑\mu (x)+{\displaystyle \frac{24^pn^{2p1}}{k^p}}{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle _{_m^n}}f(x+e_j)f(x)_X^p𝑑\mu (x)`$
$``$ $`\left({\displaystyle \frac{6^p}{3}}T^p+{\displaystyle \frac{24^pn^{2p1}}{k^p}}\right){\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle _{_m^n}}f(x+e_j)f(x)_X^p𝑑\mu (x),`$ (9)
where in (8) we used the fact that $`_j^{(k)}`$ is an averaging operator, and hence has norm $`1`$.
On the other hand
$`𝔼_\epsilon {\displaystyle _{_m^n}}f\left(x+{\displaystyle \frac{m}{2}}\epsilon \right)f(x)_X^p𝑑\mu (x)3^{p1}𝔼_\epsilon {\displaystyle _{_m^n}}𝒜^{(k)}f\left(x+{\displaystyle \frac{m}{2}}\epsilon \right)𝒜^{(k)}f(x)_X^p𝑑\mu (x)+`$ (10)
$`3^{p1}𝔼_\epsilon {\displaystyle _{_m^n}}f\left(x+{\displaystyle \frac{m}{2}}\epsilon \right)𝒜^{(k)}f\left(x+{\displaystyle \frac{m}{2}}\epsilon \right)_X^p𝑑\mu (x)+3^{p1}𝔼_\epsilon {\displaystyle _{_m^n}}𝒜^{(k)}f(x)f(x)_X^p𝑑\mu (x)`$
$``$ $`3^{p1}\left[\left({\displaystyle \frac{m}{4}}\right)^{p1}𝔼_\epsilon {\displaystyle _{_m^n}}{\displaystyle \underset{t=1}{\overset{m/4}{}}}𝒜^{(k)}f\left(x+2t\epsilon \right)𝒜^{(k)}f(x+(2t2)\epsilon )_X^pd\mu (x)+2𝔼_\epsilon {\displaystyle _{_m^n}}𝒜^{(k)}f(x)f(x)_X^p𝑑\mu (x)\right]`$
$``$ $`3^{p1}\left[\left({\displaystyle \frac{m}{4}}\right)^p𝔼_\epsilon {\displaystyle _{_m^n}}{\displaystyle \underset{t=1}{\overset{m/4}{}}}𝒜^{(k)}f\left(y+\epsilon \right)𝒜^{(k)}f(x\epsilon )_X^pd\mu (x)+2k^pn^{p1}{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle _{_m^n}}f(x+e_j)f(x)_X^p𝑑\mu (x)\right]`$
$``$ $`\left[3^{p1}\left({\displaystyle \frac{m}{4}}\right)^p\left(1+{\displaystyle \frac{1}{k}}\right)^{p(n1)}\left({\displaystyle \frac{6^p}{3}}T^p+{\displaystyle \frac{24^pn^{2p1}}{k^p}}\right)+{\displaystyle \frac{2(3kn)^p}{3n}}\right]{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle _{_m^n}}f(x+e_j)f(x)_X^p𝑑\mu (x)`$ (11)
$``$ $`5^pm^pT^p{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle _{_m^n}}f(x+e_j)f(x)_X^p𝑑\mu (x),`$ (12)
where in (10) we used Lemma 2.2, in (11) we used (8), and (12) is true if $`4n^{21/p}k\frac{3m}{2n^{11/p}}`$, which is a valid choice of $`k`$ if $`m3n^{32/p}`$. ∎
###### Remark 2.3.
If a metric space has Enflo type $`p`$ then it also has scaled Enflo type $`p`$. This follows from a straightforward modification of Lemma 2.4 in . We do not know if scaled Enflo type $`p`$ implies Enflo type $`p`$. In the category of Banach spaces, a positive answer to this question would show that Enflo type $`p`$ is equivalent to Rademacher type $`p`$, resolving positively Enflo’s problem . We do know that for Banach spaces, scaled Enflo type $`p`$ implies Enflo type $`p^{}`$ for all $`p^{}<p`$, and that scaled Enflo type and Enflo type coincide for UMD Banach spaces.
###### Remark 2.4.
The idea of scaling by $`\frac{m}{2}`$ in the definition of scaled Enflo type originates from the definition of metric cotype introduced in , which involves a similar scaling procedure. In the case on non-linear type it is possible that this scaling is not necessary, i.e. that Enflo type is equivalent to Rademacher type. However, as shown in , in the context of metric cotype the scaling is necessary\- we refer to for more details. |
warning/0506/hep-th0506011.html | ar5iv | text | # Clifford Algebras in Physics
## 1 Introduction
The real Clifford algebra $`𝒞_{t,s}`$ is the associative algebra generated by a unit $`1`$ and $`d=t+s`$ elements $`e_1,\mathrm{},e_d`$ satisfying
$`e_Me_N+e_Ne_M=2\eta _{MN},\text{with }\eta _{MN}=\{\begin{array}{cc}0\hfill & \text{if }MN\hfill \\ 1\hfill & \text{if }M=N=1,\mathrm{},t\hfill \\ 1\hfill & \text{if }M=N=t+1,\mathrm{},d.\hfill \end{array}`$ (1.4)
It is well-known that the Clifford algebra $`𝒞_{t,s}`$ can be represented by $`2^{\left[\frac{d}{2}\right]}\times 2^{\left[\frac{d}{2}\right]}`$ matrices called the Dirac matrices (with $`[a]`$ the integer part of $`a`$). Since $`𝔰𝔬(t,s)𝒞_{t,s}`$, the $`2^{\left[\frac{d}{2}\right]}`$dimensional complex vector space on which the Dirac matrices act is also a representation of $`𝔰𝔬(t,s)`$, the spinor representation. Thus on a physical ground Clifford algebras are intimately related to fermions. The four dimensional Clifford algebra, or more precisely its matrix representation, was introduced in physics by Dirac when he was looking for a relativistic first order differential equation extending the Schrödinger equation. This new equation called now, the Dirac equation, is in fact a relativistic equation describing spin one-half fermions as the electron.
Among the various applications of Clifford algebras in physics and mathematical physics, we will show how they are central tools in the construction of supersymmetric theories. In particle physics there are two types of particles: bosons and fermions. The former, as the photon, have integer spin and the latter, as the electron, have half-integer spin. Properties (mass, spin, electric charge etc.) or the way particles interact together are understood by means of Lie algebras (describing space-time and internal symmetries). Supersymmetry is a symmetry different from the previous ones in the sense that it is a symmetry which mixes bosons and fermions. Supersymmetry is not described by Lie algebras but by Lie superalgebras and contains generators in the spinor representation of the space-time symmetry group.
In this lecture, we will show, how supersymmetry is intimately related to Clifford algebras. Along all the steps of the construction, a large number of details will be given. Section 2 will be mostly mathematical and devoted to the definition, and the classification of Clifford algebras. In section 3, we will show how Clifford algebras are related to special relativity and to spinors. Section 4 is a technical section, central for the construction of supersymmetry. Some basic properties of the Dirac $`\mathrm{\Gamma }`$matrices (the matrices representing Clifford algebras) will be studied allowing to introduce different types of spinors (Majorana, Weyl). We will show how the existence of these types of spinors is crucially related to the space-time dimension. Useful details will be summarized in some tables. Finally, in section 5 we will construct explicitly supersymmetric algebras and their representations with a special emphasis to the four, ten and eleven-dimensional space-time. We will also show how lower dimensional supersymmetry are related to higher dimensional supersymmetry by dimensional reduction.
## 2 Definition and classification
In this section we give the definition of Clifford algebras. We then show that it is possible to characterize Clifford algebras as matrix algebras and show that the properties of real Clifford algebras depend on the dimension modulo $`8`$.
### 2.1 Definition
Consider $`E`$ a $`d`$dimensional vector space over the field $`𝕂`$ ($`𝕂=`$ or $`𝕂=`$). Let $`Q`$ be a quadratic form, non-degenerate of signature $`(t,s)`$ on $`E`$. The quadratic form $`Q`$ naturally defines a symmetric scalar product on $`E`$
$`\eta (x,y)={\displaystyle \frac{1}{4}}(Q(x+y)Q(xy)),\text{ for all }x,y\text{ in }E.`$ (2.1)
Denote $`\eta =\text{Diag}(\underset{t}{\underset{}{1,\mathrm{},1}},\underset{s}{\underset{}{1,\mathrm{},1}})`$ the tensor metric in an orthonormal basis. Denote also $`\eta _{MN},1M,Nd`$ the matrix elements of $`\eta `$ and $`\eta ^{MN}`$ the matrix elements of the inverse matrix $`\eta _{MN}=\eta ^{MN}=\{\begin{array}{cc}0,\hfill & MN\hfill \\ 1,\hfill & M=N=1,\mathrm{},t\hfill \\ 1,\hfill & M=N=t+1,\mathrm{}.d\hfill \end{array}`$ .
###### Definition 2.1
1) The associative $`𝕂`$algebra $`𝒞_{t,s}^𝕂`$ generated by a unit $`1`$ and $`d=t+s`$ generators $`e_M,M=1,\mathrm{},t+s`$ satisfying
$`\{e_M,e_N\}=e_Me_N+e_Ne_M=2\eta _{MN}`$ (2.2)
is called the Clifford algebra of the quadratic form $`Q`$.
2) The dimension of $`𝒞_{t,s}^𝕂`$ is $`2^d`$ and a convenient basis is given by
$`1,e_{M_1},e_{M_1}e_{M_2},\mathrm{},e_{M_1}e_{M_2}\mathrm{}e_{M_{d1}},e_1e_2\mathrm{}e_d,1M_1<M_2<\mathrm{}<M_{d1}d.`$ (2.3)
###### Remark 2.2
The case where $`Q(x)=0`$ for all $`x`$ in $`𝕂`$ will be considered latter on (see remark 4.1) and corresponds to Grassmann algebras.
###### Remark 2.3
Clifford algebras could have been defined in a more formal way as follow. Let $`T(E)=𝕂E\left(EE\right)\mathrm{}`$ be the tensor algebra over $`E`$ and let $`(Q)`$ be the two-sided ideal generated by $`xxQ(x)1,xE`$ in $`T(E)`$. The quotient algebra $`T(E)/(Q)`$ is the Clifford algebra $`𝒞_{t,s}^𝕂`$. The canonical map $`ı_Q:E𝒞_{t,s}^𝕂`$ given by the composition $`ET(E)𝒞_{t,s}^𝕂`$ is an injection.
###### Remark 2.4
If we set $`\mathrm{\Lambda }(E)=𝕂E\mathrm{\Lambda }^2(E)\mathrm{}\mathrm{\Lambda }^d(E)`$ to be the exterior algebra on $`E`$, then the canonical map $`ı_Q`$ extends to a vector-space isomorphism $`\widehat{ı}_Q:\mathrm{\Lambda }(E)𝒞_{t,s}^𝕂`$ given on $`\mathrm{\Lambda }^n(E)`$ by
$$\widehat{ı}_Q(v_1\mathrm{}v_n)=\underset{\sigma S_n}{}\frac{1}{n!}v_{\sigma (1)}\mathrm{}v_{\sigma (n)}$$
with $`S_n`$ the group of permutations with $`n`$ elements.
### 2.2 Structure
We call $`𝒞_0_{t,s}^𝕂`$ the subalgebra of $`𝒞_{t,s}^𝕂`$ generated by an even product of generators $`e_M`$ and $`𝒞_1_{t,s}^𝕂`$ the vector space (which is not a subalgebra) of $`𝒞_{t,s}^𝕂`$ generated by an odd product of generators $`e_M`$. We also introduce
$`\epsilon =e_1e_2\mathrm{}e_d`$ (2.4)
corresponding to the product of all the generators of $`𝒞_{t,s}^𝕂.`$ An easy calculation gives
$`\epsilon ^2=(1)^{\frac{d(d1)}{2}}e_1^2\mathrm{}e_M^2=(1)^{\frac{d(d1)}{2}+s}=\{\begin{array}{cc}(1)^{\frac{ts}{2}}\hfill & d\text{ even}\hfill \\ (1)^{\frac{ts1}{2}}\hfill & d\text{ odd.}\hfill \end{array}`$ (2.7)
Furthermore, when $`d`$ is even $`\epsilon `$ commutes with $`𝒞_0_{t,s}^𝕂`$ and anticommutes with $`𝒞_1_{t,s}^𝕂`$ and when $`d`$ is odd $`\epsilon `$ commutes both with $`C_0_{t,s}^𝕂`$ and $`𝒞_1_{t,s}^𝕂`$. It can be shown that the centre $`(𝒞_{t,s})=𝕂`$ (resp. $`(𝒞_{t,s})=𝕂𝕂\epsilon `$) when $`d`$ is even (resp. when $`d`$ is odd). The element $`\epsilon `$ will be useful to decide whether or not Clifford algebras are simple. We first characterize real Clifford algebras that we denote from now on $`𝒞_{t,s}`$ to simplify notation. Then, we study the case of complex Clifford algebras $`𝒞_{t,s}^{}=𝒞_{t,s}_{}`$.
#### 2.2.1 Real Clifford algebras
From (2.7) $`\epsilon ^2=1`$ for $`ts=0(\text{mod. }4)`$ when $`d`$ is even and $`\epsilon ^2=1`$ for $`ts=1(\text{mod. }4)`$ when $`d`$ is odd and thus in these cases $`P_\pm =\frac{1}{2}(1\pm \epsilon )`$ define orthogonal projection operators $`(P_+P_{}=0,P_++P_1=1)`$. These projectors will enable us to define ideals of $`𝒞_{t,s}`$ and $`𝒞_0_{t,s}`$ .
###### Proposition 2.5
Let $`𝒞_{t,s}`$ be a real Clifford algebra and let $`P_\pm `$ be defined as above.
* When $`d`$ is odd and $`ts=1(\text{mod. }4)`$ then $`P_\pm `$ belong to the centre $`(𝒞_{t,s})`$ and $`𝒞_{t,s}=P_+𝒞_{t,s}P_{}𝒞_{t,s}`$. This means that $`𝒞_{t,s}`$ is not simple ($`P_\pm 𝒞_{t,s}`$ are two ideals of $`C_{t,s})`$.
* When $`d`$ is odd, $`𝒞_0_{t,s}`$ is simple.
* When $`d`$ is even, $`𝒞_{t,s}`$ is simple.
* When $`d`$ is even and $`ts=0\text{ (mod. }4)`$ then $`P_\pm `$ belong to the centre $`(𝒞_0{}_{t,s}{}^{})`$ and thus we have $`𝒞_0{}_{t,s}{}^{}=P_+𝒞_0{}_{t,s}{}^{}P_{}𝒞_0_{t,s}`$. This means that $`𝒞_0_{t,s}`$ is not simple ($`P_\pm 𝒞_0_{t,s}`$ are two ideals of $`C_0{}_{t,s}{}^{})`$.
The proof will in fact be given in table 1, section 4.
#### 2.2.2 Complex Clifford algebras
In this case, $`E`$ is a complex vector space, and the tensor metric $`\eta `$ can always be chosen to be Euclidien, i.e. $`\eta =\mathrm{Diag}(1,\mathrm{},1)`$. We denote now $`\overline{𝒞}_d`$ a complex Clifford algebras generated by $`d`$ generators. When we are over the field of complex numbers, a projection operator can always be defined. Indeed if $`\epsilon ^2=1,(i\epsilon )^2=1`$, so we set $`P_\pm =\frac{1}{2}(1\pm \epsilon )`$ if $`\epsilon ^2=1`$ and $`P_\pm =\frac{1}{2}(1\pm i\epsilon )`$ if $`\epsilon ^2=1`$.
###### Proposition 2.6
Let $`\overline{𝒞}_d`$ be a complex Clifford algebra, we have the following structure:
* when $`d`$ is odd, $`\overline{𝒞}_d=P_+\overline{𝒞}_dP_{}\overline{𝒞}_d`$ is not simple;
* when $`d`$ is even, $`\overline{𝒞}_d`$ is simple;
* when $`d`$ is odd, $`\overline{𝒞}_0_d`$ is simple;
* when $`d`$ is even, $`\overline{𝒞}_0{}_{d}{}^{}=P_+\overline{𝒞}_0{}_{d}{}^{}P_{}\overline{𝒞}_0_d`$ is not simple.
### 2.3 Classification
We now show that it is possible to characterize all Clifford algebras as matrix algebras . (See also and references therein.) We first study the case of real Clifford algebras. The interesting point in this classification is two-fold. Firstly, the knowledge of $`𝒞_{1,0},𝒞_{0,1},𝒞_{2,0},𝒞_{0,2},𝒞_{1,1}`$ determines all the other $`𝒞_{t,s}`$. Secondly, the properties of $`𝒞_{t,s}`$ depend on $`ts`$ mod. $`8`$.
#### 2.3.1 Real Clifford algebras
In these family of algebras we may add $`𝒞_{0,0}=`$. From the definition 2.1 one can show
$`\begin{array}{ccc}𝒞_{0,0}=\hfill & 𝒞_{1,0}=\hfill & 𝒞_{0,1}=\hfill \\ 𝒞_{2,0}=_2()\hfill & 𝒞_{1,1}=_2()\hfill & 𝒞_{0,2}=\hfill \end{array}`$ (2.10)
with $``$ the quaternion algebra and $`_n(𝔽)`$ the $`n\times n`$ matrix algebra over the field $`𝔽=,`$ or $``$. Recall that the quaternion algebra is generated by three imaginary units $`i,j,k=ij`$ satisfying $`i^2=j^2=1`$ and $`ij+ji=0`$ which is precisely the definition of $`𝒞_{0,2}`$. For the algebra $`𝒞_{2,0}`$ if we set $`h_{11}=\frac{1}{2}(1+e_1),h_{22}=\frac{1}{2}(1e_1),h_{12}=h_{11}e_2`$ and $`h_{21}=h_{22}e_2`$ one can check that they are independent and that they satisfy the multiplication law of $`_2()`$. Similar simple arguments give the identifications (2.10). Next, notice the
###### Proposition 2.7
We have the following isomorphisms
$`(1)𝒞_{t,s}_{}𝒞_{2,0}𝒞_{s+2,t},(2)𝒞_{t,s}_{}𝒞_{1,1}𝒞_{t+1,s+1},(3)𝒞_{t,s}_{}𝒞_{0,2}𝒞_{s,t+2}.`$ (2.11)
Proof: To prove the first isomorphism, introduce $`\left\{e_M\right\}_{M=1,\mathrm{},d}`$ (resp. $`\{f_1,f_2\}`$) the generators of $`𝒞_{t,s}`$ (resp. $`𝒞_{2,0}`$). Then, one can show that $`\left\{g_M=e_Mf_1f_2,g_{d+1}=1f_1,g_{d+2}=1f_2\right\}_{M=1,\mathrm{},d}`$ (with $`1`$ the unit element of $`𝒞_{t,s}`$) satisfy the relations (2.2) for $`𝒞_{s+2,t}`$. To end the proof we just have to check that the $`g_{M_1}\mathrm{}g_M_{\mathrm{}},1M_1<\mathrm{}<M_{\mathrm{}}d+2,\mathrm{}=1,\mathrm{},d+2`$ are independent. A similar proof works for the two other isomorphisms. QED
Finally, we recall some well-known isomorphisms of real matrix algebras
###### Proposition 2.8
$`\begin{array}{cc}(1)_m()_n()_{mn}(),\hfill & (5)_{},\hfill \\ (2)_m()_m(),\hfill & (6)_{}_2(),\hfill \\ (3)_m()_{}_m(),\hfill & (7)_{}_4().\hfill \\ (4)_m()_{}_m(),\hfill & \end{array}`$ (2.16)
Proof: The only non-trivial isomorphisms to be proved are those involving the quaternions. Recall that $`\left\{q_0\sigma _0+q_1(i\sigma _1)+q_2(i\sigma _2)+q_3(i\sigma _3),q_0,q_1,q_2,q_3\right\}`$ with
$`\sigma _0=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right),\sigma _1=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right),\sigma _2=\left(\begin{array}{cc}0& i\\ i& 0\end{array}\right),\sigma _3=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right),`$ (2.17)
the Pauli matrices, is a faithful representation of the quaternions in $`_2()`$. Thus, $`(6)`$ is easily proven.
To prove (7), let $`q=q_01+q_1i+q_2j+q_3k,q^{}=q_0^{}1+q_1^{}i+q_2^{}j+q_3^{}k`$ be two given quaternions and set $`\overline{q}^{}=q_0^{}1q_1^{}iq_2^{}jq_3^{}k`$. Define now
$`f_{q,q^{}}:`$ $``$ $``$
$`x`$ $``$ $`x^{}=qx\overline{q}^{}.`$
It is now a matter of calculation to show that $`f_{q,q^{}}`$ can be represented by a $`4\times 4`$ real matrix acting on the component of the quaternion $`x`$. This proves $`(7)`$. QED
From the isomorphisms (2.11) and (2.16) one is able to calculate all the algebras $`𝒞_{d,0}`$ and $`𝒞_{0,d}`$. As an illustration, we give the result for $`C_{0,7}:`$
$`𝒞_{0,7}`$ $``$ $`𝒞_{5,0}𝒞_{0,2}𝒞_{0,3}𝒞_{2,0}𝒞_{0,2}𝒞_{1,0}𝒞_{0,2}𝒞_{2,0}𝒞_{0,2}`$
$``$ $`\left(\right)_2()`$
$``$ $`\left(\right)_2()`$
$``$ $`\left(_4()_4()\right)_2()`$
$``$ $`_8()_8().`$
The Clifford algebras $`𝒞_{d,0},𝒞_{0,d},0d8`$ are given in table 1.
We observe as claimed in the previous subsection that only $`𝒞_{1,0},𝒞_{0,3},𝒞_{5,0},𝒞_{0,7}`$ are not simple (see Proposition 2.5 (i)).
Next the identity
$`𝒞_{0,2}𝒞_{2,0}𝒞_{0,2}𝒞_{2,0}𝒞_{0,8}𝒞_{8,0}_{16}()`$ (2.18)
yields
$`𝒞_{d+8,0}𝒞_{d,0}𝒞_{8,0},𝒞_{0,d+8}𝒞_{0,d}𝒞_{0,8}.`$ (2.19)
Thus if $`𝒞_{d,0}_n(𝔽)`$, ($`𝔽=,`$ or $``$), then $`𝒞_{d+8,0}_{16n}(𝔽)`$. This ends the classification of the algebras $`𝒞_{d,0}`$ and $`𝒞_{0,d}`$ . The case of space-time with arbitrary signature are obtained from $`𝒞_{d,0}`$ and $`𝒞_{0,d}`$ through the identity
$`𝒞_{t,s}`$ $``$ $`𝒞_{ts,0}\underset{}{𝒞_{1,1}\mathrm{}𝒞_{1,1}},ts0`$
$`s\text{ times}`$
$``$ $`𝒞_{ts,0}_{2^s}(),t>s`$
since $`𝒞_{1,1}_2().`$ Thus if $`𝒞_{ts,0}_n(𝔽)`$, ($`𝔽=,`$ or $``$), then $`𝒞_{t,s}_{2^sn}(𝔽)`$. We summarize in table 2 the results for the Clifford algebras $`𝒞_{t,s}`$ with an arbitrary signature.
The only non-simple algebras are those when $`ts=1\text{ mod. }4`$.
#### 2.3.2 Complex Clifford algebras
Finally we give the classification for complex Clifford algebras. The complex Clifford algebras $`𝒞_{t,s}^{}𝒞_{t,s}_{}`$ are obtained by complexification of the real Clifford algebras $`𝒞_{t,s}`$. Of course the complexification of the algebra $`𝒞_{t,s}`$ depends only on $`d=t+s`$ and it will be denoted $`\overline{𝒞}_d`$. Using tables 1 and 2 we obtain
$`\overline{𝒞}_d=\{\begin{array}{cc}_{2^{\left[\frac{d}{2}\right]}}()\hfill & d\text{ even}\hfill \\ _{2^{\left[\frac{d}{2}\right]}}()_{2^{\left[\frac{d}{2}\right]}}()\hfill & d\text{ odd}.\hfill \end{array}`$ (2.23)
## 3 Clifford algebras in relation to special relativity
After recalling the definition of the basic group of invariance of special relativity, we show that the representation spaces on which Clifford algebras act correspond to spinors. We then obtain in a natural way the group $`\text{Pin}(t,s)`$ which is a non-trivial double covering of the group $`O(t,s)`$. We further identify the Lie algebra $`𝔰𝔬(t,d)𝒞_{t,s}`$.
### 3.1 Poincaré and Lorentz groups
The basic group of special relativity (i.e in a four dimensional Minkowskian space-time with a metric of signature $`(1,3)`$) can be easily extended to a $`d`$dimensional pseudo-Euclidien space of signature $`(t,s)`$. Let $`i_i,\mathrm{},i_{t+s}`$ be an orthonormal basis of the vector space $`E`$. The directions $`i_1,\mathrm{},i_t`$ are time-like ($`\eta (i_M,i_M)=1>0,1Mt`$) and the directions $`i_{t+1},\mathrm{},i_{s+t}`$ are space-like ($`\eta (i_M,i_M)=1<0,t+1Mt+s`$) . The Lorentz group is defined by $`O(t,s)=\{fGL(E):\eta (f(x),f(y))=\eta (x,y),x,yE\}`$. If $`\mathrm{\Lambda }`$ denotes a matrix representation of $`f`$, from $`f(x)=\mathrm{\Lambda }x,\eta (x,y)=x^t\eta y`$ and $`\eta (f(x),f(y))=\eta (x,y)`$, we get $`\mathrm{\Lambda }^t\eta \mathrm{\Lambda }=\eta `$ with $`\mathrm{\Lambda }^t`$ the transpose of the matrix $`\mathrm{\Lambda }`$, and $`\mathrm{\Lambda }`$ is a transformation preserving the tensor metric. The Lorentz group has four connected components
$`=O(t,s)=_+^{}_+^{}_{}^{}_{}^{}`$ (3.1)
where $`_\pm ^{},_\pm ^{}`$ represent elements of determinant $`\pm 1`$ and $``$ (resp. $``$) represents elements with positive (negative) temporal signature. Let $`R_i,i=1,\mathrm{},d`$ be the reflections in the hyperplane perpendicular to the $`i^{\text{th}}`$ direction. $`R_1,\mathrm{},R_t`$ are time-like reflections and $`R_{t+1},\mathrm{}R_{t+s}`$ are space-like reflections. More generally one can consider a reflection $`R(v)`$ orthogonal to a given direction $`vE`$ such that $`\eta (v,v)0`$ said to be time-like if $`\eta (v,v)>0`$ and space-like if $`\eta (v,v)<0`$. An element of $`O(t,s)`$ is given by a products of certain numbers of such reflections. The structure of the various components of the Lorentz group are as follow
$`\begin{array}{cc}_+^{},\hfill & \text{ is a continuous group, }\text{i.e.}\text{ is associated to some Lie algebra }\hfill \\ _{}^{}=R(v)_+^{},\hfill & \text{ where }v\text{ is a space-like direction, say }i_{t+1}\hfill \\ _{}^{}=R(v)_+^{},\hfill & \text{ where }v\text{ is a time-like direction, say }i_1\hfill \\ _+^{}=R(v)R(v^{})_+^{},\hfill & \text{ where }v\text{ is a time-like direction, say }i_1\hfill \\ & \text{ and where }v^{}\text{ is a space-like direction, say }i_{t+1}.\hfill \end{array}`$ (3.7)
In other words, if $`R(v_1)\mathrm{}R(v_n)`$ is a product of (i) an even number of space-like and time-like reflections it belongs to $`_+^{}`$, (ii) an odd number of space-like and an even number of time-like reflections it belongs to $`_{}^{}`$, (iii) an even number of space-like and an odd number of time-like reflections it belongs to $`_{}^{}`$, (iv) an odd number of space-like and time-like reflections to it belongs to $`_+^{}`$. When the signature is $`(1,d1)`$, $`R(e_1)=T`$ is the operator of time reversal (it changes the direction of the time), and when in addition $`d`$ is even $`R(e_2)\mathrm{}R(e_d)=P`$ is the parity operator (it corresponds to a reflection with respect to the origin). Furthermore, one can identify several subgroups of $`O(t,s)`$:
$`O(t,s)=,SO(t,s)=_+^{}_+^{},SO_+(t,s)=_+^{}.`$ (3.8)
$`SO(t,s)`$ is constituted of an even product of reflections and $`SO_+(t,s)`$ an even product of time-like and space-like reflections. Note that with an Euclidien metric, $`O(d)`$ has only two connected components. Moreover, none of these groups are simple connected.
For further use, we introduce the generators of $`𝔰𝔬(t,s)`$ the Lie algebra of $`SO_+(t,s)`$. A conventional basic is given by the $`L_{MN}=L_{NM}`$ which corresponds to the generators which generate the “rotations” in the plane $`(MN)`$. We also introduce $`P_M`$ ($`1Md`$) the generators of the space-time translations.<sup>1</sup><sup>1</sup>1If $`x^M,M=1,\mathrm{},d`$ denote the components in the $`d`$dimensional vector space $`E`$ and we set $`x_N=\eta _{NM}x^M`$, we have $`L_{MN}=(x_M\frac{}{x^N}x_N\frac{}{x^M})`$ and $`P_M=\frac{}{x^M}`$. The generators $`L_{MN}`$ and $`P_M`$ generate the so-called Poincaré algebra or the inhomogeneous Lorentz algebra noted $`𝔦𝔰𝔬(t,s)`$ and satisfy
$`[L_{MN},L_{PQ}]`$ $`=`$ $`\eta _{MP}L_{NQ}+\eta _{NP}L_{MQ}\eta _{MQ}L_{PN}+\eta _{NQ}L_{PM},`$
$`[L_{MN},P_P]`$ $`=`$ $`\eta _{NP}P_M\eta _{MP}P_N.`$ (3.9)
### 3.2 Universal covering group of the Lorentz group
In this section we only consider real Clifford algebras $`𝒞_{t,s}`$. We furthermore identify a vector $`v=x^Mi_ME`$ with an element $`x^Me_M𝒞_{t,s}`$ (see remark 2.3). With such an identification, we have $`x^2=((x^1)^2+\mathrm{}+(x^t)^2(x^{t+1})^2\mathrm{}(x^{t+s})^2=\eta (x,x).`$
###### Definition 3.1
The Clifford group $`\mathrm{\Gamma }_{t,s}`$ is the subset of invertible elements $`x`$ of $`𝒞_{t,s}`$ such that $`vE,xvx^1E`$.
It is clear that invertible elements of $`E`$ belong to $`\mathrm{\Gamma }_{t,s}`$ (this excludes the null vectors i.e the vectors such that $`\eta (v,v)=0`$). If we take $`xE𝒞_{t,s}`$ invertible, since $`xv+vx=2\eta (x,v)`$, we have
$`xvx^1={\displaystyle \frac{1}{x^2}}xvx=\left(v{\displaystyle \frac{2\eta (x,v)}{x^2}}x\right)=R(x)(v),`$ (3.10)
which corresponds to a symmetry in the hyperplane perpendicular to $`x`$. More generally, for $`s\mathrm{\Gamma }_{t,s}`$, the transformation $`\rho (s)`$ defined by $`\rho (s)(x)=sxs^1`$ belong to $`O(t,s)`$. However, the representation $`\rho :\mathrm{\Gamma }_{t,s}GL(E)`$ is not faithful because if $`x\mathrm{\Gamma }_{t,s}`$, then $`ax\mathrm{\Gamma }_{t,s}`$ $`(a^{})`$ and $`\rho (ax)=\rho (x)`$. A standard way to distinguish the elements $`x`$ and $`ax`$ is to introduce a normalisation. For that purpose we define $`\overline{\text{ }}:𝒞_{t,s}𝒞_{t,s}`$ by
$`\overline{e_M}=(1)^{t+1}e_M,\overline{e_{M_1}.\mathrm{}.e_{M_{k1}}.e_{M_k}}=\overline{e_{M_k}}.\overline{e_{M_{k1}}}.\mathrm{}.\overline{e_{M_1}},`$ (3.11)
(see remark 4.8). Next, we define $`N(x)=x\overline{x}`$. Note that in the case of the quaternions ($`t=0,s=2`$), we have $`\overline{i}=i,\overline{j}=j,\overline{k}=k`$. Now, with this definition of the norm, for $`xE`$ we have
$`(1)^{t+1}N(x)=x^2\{\begin{array}{cc}>0\hfill & \text{ if }x\text{ is time-like }\hfill \\ <0\hfill & \text{ if }x\text{ is space-like }\hfill \end{array}`$ (3.14)
###### Proposition 3.2
(Proposition 3.8 of ) If $`x\mathrm{\Gamma }_{t,s}`$ then $`N(x)^{}`$.
As a consequence for $`x,y\mathrm{\Gamma }_{t,s}`$ we have $`N(xy)=N(x)N(y)`$. Then, the “norm” $`N`$ enables us to define definite subgroups of the Clifford group $`\mathrm{\Gamma }_{t,s}`$. The first subgroup which can be defined is
$`\text{Pin}(t,s)`$ $`=`$ $`\left\{x\mathrm{\Gamma }_{t,s}\text{ s.t. }|N(x)|=1\right\}.`$ (3.15)
By construction it is easy to see that $`\rho (\mathrm{\Gamma }_{t,s})=\rho (\text{Pin}(t,s))`$ and that $`xE`$ invertible $`x\text{Pin}(t,s)`$. Moreover, since an element of $`O(t,s)`$ is given by a product of a given number of reflections, we have $`\rho (\text{Pin}(t,s))O(t,s)`$. Conversely it has been shown (, proposition 3.10) that $`\rho (\text{Pin}(t,s))O(t,s)`$, thus $`\rho (\text{Pin}(t,s))=O(t,s)`$. A generic element of $`\text{Pin}(t,s)`$ is then given by
$`s=v_1v_2\mathrm{}v_n`$ (3.16)
with $`v_i,i=1,\mathrm{},n`$ invertible elements of $`E`$ (thus $`\rho (s)`$ corresponds to the transformation given by $`R(v_1)R(v_2)\mathrm{}R(v_n)`$). Assume now, that there are $`p`$ time-like vectors and $`q`$ space-like vectors in $`s`$ ($`p+q=n`$), then $`N(s)=(1)^{p(t+1)+qt}=(1)^{nt+p}.`$
1. If $`n`$ is even then $`\rho (s)SO(t,s)`$.
2. If $`n`$ is even and $`N(s)>0`$ then $`\rho (s)SO_+(t,s)`$.
3. If $`n`$ is even and $`N(s)<0`$ then $`\rho (s)_+^{}`$. For instance when $`d`$ is an even number and $`t=1,s=d1`$, $`\rho (\epsilon )=\rho (e_1\mathrm{}e_d)_+^{}`$ and corresponds to the $`PT`$ inversion ($`P`$ being the operator of parity transformation and $`T`$ the operator of time reversal.)
4. If $`n`$ is odd and $`N(s)(1)^{t+1}>0`$ then $`\rho (s)_{}^{}`$. For instance when $`(t,s)=(1,d1),\rho (e_1)_{}^{}`$ and corresponds to the operator of time reversal.
5. If $`n`$ is odd and $`N(s)(1)^{t+1}<0`$ then $`\rho (s)_{}^{}`$. For instance when $`d`$ is even and $`(t,s)=(1,d1),\rho (e_1\mathrm{}e_d))_{}^{}`$ and corresponds to the operator of parity transformation.
We can now define the various subgroup of $`\mathrm{\Gamma }_{t,s}`$:
$`\text{Pin}(t,s)`$ $`=`$ $`\left\{x\mathrm{\Gamma }_{t,s}\text{ s.t. }|N(x)|=1\right\}`$
$`\text{Spin}(t,s)`$ $`=`$ $`\text{Pin}(t,s)𝒞_0_{t,s}`$ (3.17)
$`\text{Spin}_+(t,s)`$ $`=`$ $`\left\{x\text{Spin}(t,s)\text{ s.t. }N(x)=1\right\}`$
### 3.3 The Lie algebra of $`\text{Spin}_+(t,s)`$
Among the various groups of the previous subsection, only $`\text{Spin}_+(t,s)`$ is a connected Lie group. If one introduce the $`d(d1)/2`$ elements
$`S_{MN}={\displaystyle \frac{1}{4}}(e_Me_Ne_Ne_M),1MNd`$ (3.18)
a direct calculation gives
$`[S_{MN},e_P]`$ $`=`$ $`\eta _{NP}e_M\eta _{MP}e_N`$ (3.19)
$`[S_{MN},S_{PQ}]`$ $`=`$ $`\eta _{MP}S_{NQ}+\eta _{NP}S_{MQ}\eta _{MQ}S_{PN}+\eta _{NQ}S_{PM}.`$
This means that $`S_{MN}`$ generate the Lie algebra $`𝔰𝔬(t,s)`$ ($`𝔰𝔬(t,s)𝒞_{t,s}`$) and that $`e_M`$ are in the vector representation of $`𝔰𝔬(t,s)`$. Furthermore, if we define
$`e_{M_1\mathrm{}M_{\mathrm{}}}^{(\mathrm{})}={\displaystyle \frac{1}{\mathrm{}!}}{\displaystyle \underset{\sigma \mathrm{\Sigma }_{\mathrm{}}}{}}ϵ(\sigma )e_{M_{\sigma (1)}}\mathrm{}e_{M_{\sigma (\mathrm{})}},`$ (3.20)
(with $`ϵ(\sigma )`$ the signature of the permutation $`\sigma `$) using (3.19) one can show that they are in the $`\mathrm{}^{\text{ th}}`$antisymmetric representation of $`𝔰𝔬(t,s)`$ (see remark 2.4).
###### Proposition 3.3
The group $`\text{Spin}_+(t,s)`$ is a non-trivial double covering group of $`SO_+(t,s).`$
Proof: First notice that if $`x\text{Spin}_+(t,s)`$ then $`x\text{Spin}_+(t,s)`$. Then, we show that there exist a continuous path in $`\text{Spin}_+(t,s)`$ which connects $`1`$ to $`1`$. Let $`i_M,i_N`$ be two space-like directions. The path $`R(\theta )=e^{\frac{1}{2}\theta e_Me_N}=\mathrm{cos}\frac{\theta }{2}+\mathrm{sin}\frac{\theta }{2}e_Me_N`$ with $`\theta [0,2\pi ]`$ is such that $`R(0)=1`$ and $`R(2\pi )=1`$ and thus connects $`1`$ to $`1`$ in $`\text{Spin}_+(t,s)`$. Which ends the proof. QED
In the same way, $`\text{Spin}(t,s)`$ is a non-trivial double covering group of $`SO(t,s)`$ and $`\text{ Pin}(t,s)`$<sup>2</sup><sup>2</sup>2This is a joke of J.- P. Serre. a non-trivial double covering group of $`O(t,s)`$.
## 4 The Dirac $`\mathrm{\Gamma }`$matrices
In this section, because we are mostly interested in Clifford algebras in relation to space-time physics, we will focus on vector spaces with a Lorentzian signature $`(1,d1)`$ or an Euclidian signature $`(0,d)`$. From now on, in the case of the $`(1,d1)`$ signature, we use the notations commonly used in the literature. The indices of the space-time components run from $`0`$ to $`d1`$, greek indices $`\mu ,\nu ,\mathrm{}=0,\mathrm{},d1`$ are space-time indices and latin indices $`i,j,\mathrm{}=1,\mathrm{},d1`$ are space indices. This section is devoted to the study of the matrix representation of Clifford algebras. The precise structure of these matrices gives rise to the type of spinors one is able to construct in a given space-time dimension: Majorana or Weyl. Properties of spinors are also studied. This section is technical but is central for the construction of supersymmetric theories in various dimensions.
### 4.1 Dirac spinors
For physical applications, we need to have matrix representations of Clifford algebras. As we have seen in section 2.3, the case where $`d`$ is even is very different from the case where $`d`$ is odd. Indeed, complex Clifford algebras are simple when $`d`$ is even and are not simple when $`d`$ is odd. This means that the representation is faithful when $`d`$ is even and not faithful when $`d`$ is odd. When $`d`$ is odd, since $`\epsilon =e_1\mathrm{}e_d`$ allows to define the two ideals of the Clifford algebra, $`\epsilon `$ will be represented by a number. Moreover, one of the ideals will be the kernel of the representation.
It is well-known that the basic building block of the matrix representation of $`\overline{C}_d`$ are the Pauli matrices (2.17) . Define the ($`2k+1`$) matrices by the tensor products of $`k`$ Pauli matrices (we thus construct $`2^k\times 2^k`$ matrices):
$`\begin{array}{cc}\mathrm{\Sigma }_1^{(k)}=\sigma _1\underset{k1}{\underset{}{\sigma _0\mathrm{}\sigma _0}},\hfill & \mathrm{\Sigma }_2^{(k)}=\sigma _2\underset{k1}{\underset{}{\sigma _0\mathrm{}\sigma _0}},\hfill \\ \mathrm{\Sigma }_3^{(k)}=\sigma _3\sigma _1\underset{k2}{\underset{}{\sigma _0\mathrm{}\sigma _0}},\hfill & \mathrm{\Sigma }_4^{(k)}=\sigma _3\sigma _2\underset{k2}{\underset{}{\sigma _0\mathrm{}\sigma _0}},\hfill \\ \mathrm{}\hfill & \mathrm{}\hfill \\ \mathrm{\Sigma }_{2\mathrm{}1}^{(k)}=\underset{\mathrm{}1}{\underset{}{\sigma _3\mathrm{}\sigma _3}}\sigma _1\underset{k\mathrm{}}{\underset{}{\sigma _0\mathrm{}\sigma _0}},\hfill & \mathrm{\Sigma }_2\mathrm{}^{(k)}=\underset{\mathrm{}1}{\underset{}{\sigma _3\mathrm{}\sigma _3}}\sigma _2\underset{k\mathrm{}}{\underset{}{\sigma _0\mathrm{}\sigma _0}},\hfill \\ \mathrm{}\hfill & \mathrm{}\hfill \\ \mathrm{\Sigma }_{2k1}^{(k)}=\underset{k1}{\underset{}{\sigma _3\mathrm{}\sigma _3}}\sigma _1,\hfill & \mathrm{\Sigma }_{2k}^{(k)}=\underset{k1}{\underset{}{\sigma _3\mathrm{}\sigma _3}},\sigma _2,\hfill \end{array}`$ (4.7)
$`\mathrm{\Sigma }_{2k+1}^{(k)}=\underset{k}{\underset{}{\sigma _3\mathrm{}\sigma _3}}.`$ (4.8)
Observing that
$`\mathrm{\Sigma }_M^{(k+1)}=\mathrm{\Sigma }_M^{(k)}\sigma _0,M=1,\mathrm{},2k,\mathrm{\Sigma }_{2k+i}^{(k+1)}=\mathrm{\Sigma }_{2k+1}^{(k)}\sigma _i,i=1,2,3`$ (4.9)
a simple recurrence on $`k`$ shows that the $`\mathrm{\Sigma }`$matrices satisfy $`\mathrm{\Sigma }_M^{(k)}\mathrm{\Sigma }_N^{(k)}+\mathrm{\Sigma }_N^{(k)}\mathrm{\Sigma }_M^{k)}=2\delta _{MN}`$. Thus
$`\gamma :`$ $`\overline{𝒞}_d`$ $`_{2^k}()`$ (4.10)
$`e_M`$ $`\mathrm{\Sigma }_M`$
($`d=2k`$ or $`d=2k+1`$) is a representation of $`\overline{C}_d`$. This representation is faithful when $`d`$ is even and non-faithful when $`d`$ is odd ($`\gamma (\epsilon )=\mathrm{\Sigma }_1^{(k)}\mathrm{\Sigma }_2^{(k)}\mathrm{}\mathrm{\Sigma }_{2k+1}^{(k)}=i^k.`$) Since the group $`\text{Spin}_+(t,s)`$ is the double covering of the group $`SO_+(t,s)`$, the representation on which the $`\mathrm{\Gamma }`$matrices act is a representation of $`\text{Spin}_+(t,s)`$. The elements of the representation space $`\overline{𝒞}_d`$ are called Dirac spinors. A Dirac spinor exists in any dimension $`d`$ and has $`2^{\left[\frac{d}{2}\right]}`$ complex components. We denote by $`\mathrm{\Psi }_D`$ a Dirac spinor.
Now having represented the algebra $`\overline{𝒞}_d`$, ($`d=2k`$ or $`d=2k+1`$) if we set $`\mathrm{\Gamma }_M=\mathrm{\Sigma }_M^{(k)},1Mt`$ and $`\mathrm{\Gamma }_M=i\mathrm{\Sigma }_M^{(k)},t+1Mt+s`$ we have a representation of the algebra $`𝒞_{t,s}`$ corresponding to a real form of $`\overline{𝒞}_d`$. In particular for a Minkowskian space-time we introduce
$`\mathrm{\Gamma }_0=\mathrm{\Sigma }_1^{(k)},\mathrm{\Gamma }_j=i\mathrm{\Sigma }_{j+1}^{(k)},j=1,\mathrm{},d1.`$ (4.11)
We also denote by
$`\mathrm{\Gamma }_{\mu \nu }={\displaystyle \frac{1}{4}}(\mathrm{\Gamma }_\mu \mathrm{\Gamma }_\nu \mathrm{\Gamma }_\nu \mathrm{\Gamma }_\mu )`$ (4.12)
the generators of the Lie algebra $`𝔰𝔬(1,d1)`$.
###### Remark 4.1
Without using the result of the section 2.3 we show that $`\overline{𝒞}_{2n}_{2^n}()`$ as follows. Set
$`a_i={\displaystyle \frac{1}{2}}\left(e_{2i}+ie_{2i+1}\right),b_i={\displaystyle \frac{1}{2}}\left(e_{2i}ie_{2i+1}\right),i=1,\mathrm{},n`$ (4.13)
which satisfy
$`a_ia_j+a_ja_i=0,b_ib_j+b_jb_i=0,a_ib_j+b_ja_i=\delta _{ij}.`$ (4.14)
Thus the $`a_i`$ and $`b_i`$ generate the fermionic oscillator algebra (i.e. the algebra which underlines the fermionic fields after quantization). Note that the $`a_i`$ alone generate the Grassmann algebra of dimension $`n`$. To obtain a representation of the algebra (4.14), we introduce the Clifford vacuum $`\mathrm{\Omega }=|\frac{1}{2},\frac{1}{2},\mathrm{},\frac{1}{2}`$ such that $`a_i\mathrm{\Omega }=0,i=1,\mathrm{},n`$. Then we obtain the representation of the algebra (4.14) by acting in all possible ways with $`b_i`$ at most once each:
$`|s_1,s_2,\mathrm{},s_n=(b_1)^{s_1+1/2}\mathrm{}(b_n)^{s_n+1/2}|{\displaystyle \frac{1}{2}},{\displaystyle \frac{1}{2}},\mathrm{},{\displaystyle \frac{1}{2}},s_1,s_2,\mathrm{},s_n={\displaystyle \frac{1}{2}},{\displaystyle \frac{1}{2}}`$ (4.15)
and thus we obtain a $`2^n`$dimensional representation. The notation with $`s_i=\frac{1}{2},\frac{1}{2}`$ (instead of $`s_i=0,1`$) seems to be unnatural, but is appropriate to identify the weight of $`|s_1,s_2,\mathrm{},s_n`$ with respect to the Cartan subalgebra of $`𝔰𝔬(2n,)`$. Furthermore, it can be shown that $`a_i^{}=b_i`$ and that the representation is unitary for the real algebra $`𝒞_{2n,0}`$. In this case, the notation reflects the property of the Dirac spinor with respect to the group $`\text{Spin}(2n)`$. A basis of the Cartan subalgebra of $`𝔰𝔬(2n)`$ can be taken to be $`\mathrm{\Gamma }_{2i2i+1}=ia_ia_i^{}\frac{i}{2}`$ and the vector $`|s_1,s_2,\mathrm{},s_n`$ is an eigenvector of $`\mathrm{\Gamma }_{2i2i+1}`$ with eigenvalue $`is_i`$. The half-integer weights show that we have a spinor representation of $`\text{Spin}(2n)`$. (The $`i`$ factor comes from the fact that $`SO(2n)`$ is the real form of $`SO(2n,)`$ corresponding to the maximal compact algebra, for which the generators are antihermitian $``$note that, in the physical literature there is no $`i`$ factor since $`\mathrm{\Gamma }_{\mu \nu }=i\frac{1}{2}\mathrm{\Gamma }_\mu \mathrm{\Gamma }_\nu `$ instead of (4.12).$``$ ) To end this remark we just notice that the fermionic generators $`(a_i,b_i)`$ are in the vector representation of $`\text{Spin}(2N)`$.
### 4.2 Majorana and Weyl spinors
In the previous subsection we have introduced Dirac spinors. In this subsection, we will see that further spinors can be defined such as Majorana spinors or Weyl spinors etc. The Weyl spinors are just a consequence of the reducibility of the (Dirac) spinors representation of $`𝔰𝔬(2n,)`$ (and of course of all its real forms). They exist in any even space-time dimensions and for any signature. If there exists a real matrix representation of the Clifford algebra $`𝒞_{t,s}`$ then one is able to consider real (or Majorana) spinors. The existence of Majorana spinors depends on the space-time dimension and of the signature of the metric.
If we denote $`𝒮`$ the vector space corresponding to the Dirac spinors, recall the Lorentz generators writes (3.18)
$`\mathrm{\Gamma }_{\mu \nu }={\displaystyle \frac{1}{4}}(\mathrm{\Gamma }_\mu \mathrm{\Gamma }_\nu \mathrm{\Gamma }_\nu \mathrm{\Gamma }_\mu ).`$ (4.16)
The matrices $`\mathrm{\Gamma }_\mu `$ act on the representation $`𝒮`$ and the representation of $`𝒞_{t,s}`$ is just $`\text{End}(𝒮)`$, the set of endomorphisms of $`𝒮`$. In the same way the matrices $`\mathrm{\Gamma }^t`$ (with <sup>t</sup> the transpose operation) act on the dual representation $`𝒮^{}`$, the matrices $`\mathrm{\Gamma }^{}`$ (with the complex conjugation<sup>3</sup><sup>3</sup>3In the mathematical literature the complex conjugate of $`\mathrm{\Gamma }_\mu `$ is denoted $`\overline{\mathrm{\Gamma }}_\mu `$.) act on the complex-conjugate representation $`\overline{𝒮}`$ and the matrices $`\mathrm{\Gamma }^{}`$ (with the hermitian conjugation) act on the representation $`\overline{𝒮}^{}`$. These representations are in fact equivalent, since we can find elements of $`\text{End}(𝒮)`$ such that
$`A\mathrm{\Gamma }_{\mu \nu }A^1=\mathrm{\Gamma }_{\mu \nu }^{},B\mathrm{\Gamma }_{\mu \nu }B^1=\mathrm{\Gamma }_{\mu \nu }^{},C\mathrm{\Gamma }_{\mu \nu }C^1=\mathrm{\Gamma }_{\mu \nu }^t`$ (4.17)
(see below). The operators $`A,B,C`$ are intertwining operators ($`B`$ intertwines the representations $`𝒮`$ and $`\overline{𝒮}`$). This is in fact related to (4.16) and
$`A\mathrm{\Gamma }_\mu A^1=\eta _A\mathrm{\Gamma }_\mu ^{},B\mathrm{\Gamma }_\mu B^1=\eta _B\mathrm{\Gamma }_\mu ^{},C\mathrm{\Gamma }_\mu C^1=\eta _C\mathrm{\Gamma }_{\mu \nu }^t,`$ (4.18)
with $`\eta _A,\eta _B`$ and $`\eta _C`$ signs which depend on $`d`$, as we will see in subsection 4.2.2.
As a direct consequence of (4.17), if $`\mathrm{\Psi }_D𝒮`$ transforms like $`\mathrm{\Psi }_D^{}=S(\alpha )\mathrm{\Psi }_D=e^{\frac{1}{2}\alpha ^{\mu \nu }\mathrm{\Gamma }_{\mu \nu }}\mathrm{\Psi }_D`$ under a Lorentz transformation ($`\alpha _{\mu \nu }`$ are the parameters of the transformation) then $`\overline{\mathrm{\Xi }}_D=\mathrm{\Xi }_D^{}A`$ and $`\mathrm{\Xi }^tC`$ belong to $`𝒮^{}`$ (i.e transform with $`S(\alpha )^1`$). Thus, $`\overline{\mathrm{\Xi }}_D\mathrm{\Psi }_D`$ and $`\mathrm{\Xi }_D^tC\mathrm{\Psi }_D`$ define invariants.
#### 4.2.1 Weyl spinors
Consider first the case of the complex Clifford algebra $`\overline{𝒞}_d`$. When $`d`$ is even the spinor representation $`𝒮`$ is reducible $`𝒮=𝒮_+𝒮_{}`$. Indeed, if we define the chirality matrix
$`\chi =i^{\left[\frac{d}{2}\right]+1}\mathrm{\Gamma }_0\mathrm{\Gamma }_1\mathrm{}\mathrm{\Gamma }_{d1}`$ (4.19)
it satisfies
$`\chi ^2=1,\{\mathrm{\Gamma }_\mu ,\chi \}=0,[\chi ,\mathrm{\Gamma }_{\mu \nu }]=0.`$ (4.20)
Hence, $`\chi `$ allows to define the complex left- and right-handed Weyl spinors. These spinors correspond to the two irreducible representations of $`\overline{𝒞}_d`$:
$`\mathrm{\Psi }_L={\displaystyle \frac{1}{2}}(1\chi )\mathrm{\Psi }_D,\mathrm{\Psi }_R={\displaystyle \frac{1}{2}}(1+\chi )\mathrm{\Psi }_D,`$ (4.21)
$`\mathrm{\Psi }_L𝒮_{},\mathrm{\Psi }_R𝒮_+`$ and $`\mathrm{\Psi }_D𝒮`$ or $`𝒮_\pm =\left\{\mathrm{\Psi }𝒮\text{ s.t. }\chi \mathrm{\Psi }=\pm \mathrm{\Psi }\right\}.`$ This means that the spaces $`𝒮_\pm `$ carry an irreducible representation of the complex algebra $`\overline{𝒞}_{0d}`$. Now, if we concider the real form $`𝒞_0_{t,s}`$ of $`\overline{𝒞}_{0d}`$, the spaces $`𝒮_\pm `$ become irreducible representations of $`𝒞_0_{t,s}`$. The generators of the two Weyl spinors of $`𝔰𝔬(1,d1)`$ are $`\mathrm{\Sigma }_{\mu \nu }^\pm =\frac{1}{2}\left(1\pm \chi \right)\mathrm{\Gamma }_{\mu \nu }`$. The Weyl spinors are called the semi-spinors in the mathematical literature.
###### Remark 4.2
As we have seen (see Proposition 2.6 (i)), the algebra $`\overline{𝒞}_0_d`$ is not simple. Then it falls into two simple ideals $`𝒞_0{}_{d}{}^{}=P_+𝒞_0{}_{d}{}^{}+P_{}𝒞_0_d`$. Moreover, the Lie algebra $`𝔰𝔬(t,s)𝒞_0_d`$, this means that a (complex) Dirac spinor is reducible into two (complex) Weyl spinors.
#### 4.2.2 Majorana spinors
Majorana spinors are real spinors. As we now see, the existence of Majorana spinors crucially depends on the dimension $`d`$ and on the metric signature. We consider the case of Lorentzian signatures $`(1,d1)`$. The general case can be easily deduced. However, for a general study see , .
The key observation in this subsection is the simple fact that (i) the Pauli matrices are hermitian (ii) the matrices $`\sigma _1,\sigma _3`$ are real and symmetric and (iii) the matrix $`\sigma _2`$ is purely imaginary and antisymmetric. Thus from (4.7) we see (take $`d=2n+1`$ odd)
$`\begin{array}{cccc}\mathrm{\Gamma }_0^{}=\mathrm{\Gamma }_0,\hfill & \mathrm{\Gamma }_0^{}=\mathrm{\Gamma }_0,\hfill & \mathrm{\Gamma }_0^t=\mathrm{\Gamma }_0,\hfill & \\ \mathrm{\Gamma }_{2i}^{}=\mathrm{\Gamma }_{2i},\hfill & \mathrm{\Gamma }_{2i}^{}=\mathrm{\Gamma }_{2i},\hfill & \mathrm{\Gamma }_{2i}^t=\mathrm{\Gamma }_{2i},\hfill & i=1,2,\mathrm{},n\hfill \\ \mathrm{\Gamma }_{2i1}^{}=\mathrm{\Gamma }_{2i1},\hfill & \mathrm{\Gamma }_{2i1}^{}=\mathrm{\Gamma }_{2i1},\hfill & \mathrm{\Gamma }_{2i1}^t=\mathrm{\Gamma }_{2i1},\hfill & i=0,2,\mathrm{},n.\hfill \end{array}`$ (4.25)
The results that we will establish here seems to depend of the choise of basis we have chosen, but in fact they are independent of this choice (see ).
Space-time of even dimension $`d=2n`$
When $`d`$ is even, we take the first $`d`$ matrices above. In this case, we have seen that complex Clifford algebras are isomorphic to some matrix algebras. This means that if we have another representation of the Clifford algebra $`\mathrm{\Gamma }_\mu ^{}(\{\mathrm{\Gamma }_\mu ^{},\mathrm{\Gamma }_\nu ^{}\}=2\eta _{\mu \nu }`$), there exists an invertible matrix $`U`$ of $`\text{End}(𝒮)`$ such that $`\mathrm{\Gamma }_\mu ^{}=U\mathrm{\Gamma }_\mu U^1`$. In particular, we have (see (4.25) and (4.17))
$`\begin{array}{cc}\chi \mathrm{\Gamma }_\mu \chi ^1=\mathrm{\Gamma }_\mu ,\hfill & \chi =i^{n+1}\mathrm{\Gamma }_0\mathrm{\Gamma }_1\mathrm{}\mathrm{\Gamma }_{2n1}\hfill \\ A\mathrm{\Gamma }_\mu A^1=\mathrm{\Gamma }_\mu ^{},\hfill & A=\mathrm{\Gamma }_0\hfill \\ B_1\mathrm{\Gamma }_\mu B_1^1=(1)^{n+1}\mathrm{\Gamma }_\mu ^{},\hfill & B_1=\mathrm{\Gamma }_2\mathrm{\Gamma }_4\mathrm{}\mathrm{\Gamma }_{2n2}\hfill \\ B_2\mathrm{\Gamma }_\mu B_2^1=(1)^n\mathrm{\Gamma }_\mu ^{},\hfill & B_2=\mathrm{\Gamma }_0\mathrm{\Gamma }_1\mathrm{}\mathrm{\Gamma }_{2n1}\hfill \\ C_1\mathrm{\Gamma }_\mu C_1^1=(1)^{n+1}\mathrm{\Gamma }_\mu ^t,\hfill & C_1=\mathrm{\Gamma }_0\mathrm{\Gamma }_2\mathrm{}\mathrm{\Gamma }_{2n2}\hfill \\ C_2\mathrm{\Gamma }_\mu C_2^1=(1)^n\mathrm{\Gamma }_\mu ^t,\hfill & C_2=\mathrm{\Gamma }_1\mathrm{\Gamma }_3\mathrm{}\mathrm{\Gamma }_{2n1}\hfill \end{array}`$ (4.32)
$`\mathrm{\Gamma }_0`$ is the only hermitian matrix, $`B_1`$ is the product of all purely imaginary matrices, $`B_2`$ is the product of real matrices, $`C_1`$ is the product of symmetric matrices and $`C_2`$ of antisymmetric matrices. Note also the relation between the matrices $`A,B,C`$: $`C_i=AB_i,i=1,2`$. We set $`\eta _1=(1)^{n+1},\eta _2=(1)^n`$. From the definition of $`B`$, we have
$`B_1B_1^{}=(1)^{\frac{(n1)(n2)}{2}}=ϵ_1,B_2B_2^{}=(1)^{\frac{n(n1)}{2}}=ϵ_2.`$ (4.33)
Now for further use, we collect the following signs
$`\begin{array}{ccccc}d=2\text{ mod. }8\hfill & ϵ_1=+,\hfill & \eta _1=+\hfill & ϵ_2=+,\hfill & \eta _2=\hfill \\ d=4\text{ mod. }8\hfill & ϵ_1=+,\hfill & \eta _1=\hfill & ϵ_2=,\hfill & \eta _2=+\hfill \\ d=6\text{ mod. }8\hfill & ϵ_1=,\hfill & \eta _1=+\hfill & ϵ_2=,\hfill & \eta _2=\hfill \\ d=8\text{ mod. }8\hfill & ϵ_1=,\hfill & \eta _1=\hfill & ϵ_2=+,\hfill & \eta _2=+.\hfill \end{array}`$ (4.38)
Majorana spinors, $`d`$ even
From (4.32), we get $`B\mathrm{\Gamma }_{\mu \nu }B^1=\mathrm{\Gamma }_{\mu \nu }^{},B=B_1,B_2,`$ so the Dirac spinor $`\mathrm{\Psi }_D`$ and $`B^1\mathrm{\Psi }_D^{}`$ transform in the same way under the group $`\text{Pin}(1,d1)`$. A Majorana spinor is a spinor that we impose to be a real spinor. It satisfies
$`\mathrm{\Psi }_M^{}=B\mathrm{\Psi }_M,B=B_1,B_2.`$ (4.39)
But taking the complex conjugate of the above equation gives $`\mathrm{\Psi }_M=B^{}\mathrm{\Psi }_M^{}=B^{}B\mathrm{\Psi }_M`$. Thus this is possible only if $`BB^{}=1`$ or when $`ϵ_1=1`$ and/or $`ϵ_2=1`$ and from (4.38) when $`d=2,4,8\text{ mod. }8`$. More precisely, looking to (4.32), we observe that the $`\mathrm{\Gamma }`$matrices can be taken to be purely real if $`\eta _1=ϵ_1=1`$ or $`\eta _2=ϵ_2=1`$ i.e. if $`d=2,8\text{ mod. }8`$ and the Dirac matrices can be taken purely imaginary if $`ϵ_1=1,\eta _1=1`$ or $`ϵ_2=1,\eta _2=1`$ that is when $`d=2,4\text{ mod. }8`$ (In this case the matrices $`\mathrm{\Gamma }_\mu `$ are purely imaginary, the matrices $`\mathrm{\Gamma }_{\mu \nu }^{(2)}`$ are real, the matrices $`\mathrm{\Gamma }_{\mu \nu \rho }^{(3)}`$ are purely imaginary etc.). The first type of real spinors will be called Majorana spinors although the second type of spinors pseudo-Majorana spinors. This result could have been deduced from table 2. Indeed, we have shown that $`𝒞_{t,s}`$ is real when $`ts=0,1,2\text{ mod. }8`$. Observing that $`𝒞_{t,s}`$ is related to $`𝒞_{s,t}`$ by the transformation $`e_Mie_M`$, we get that $`\gamma (e_M)=\mathrm{\Gamma }_M`$ are purely imaginary when $`st=1,2\text{ mod. }8`$. Thus we have:
###### Proposition 4.3
Assume $`d`$ is even.
* Majorana spinors of $`\text{P}in(1,d1)`$ exist when $`d=2,8\text{ mod. }8`$.
* Pseudo-Majorana spinors of $`\text{P}in(1,d1)`$ exist when $`d=2,4\text{ mod. }8`$.
$`SU(2)`$Majorana spinors, $`d`$ even
As we have seen, a Majorana spinor is a spinor which satisfies $`\mathrm{\Psi }_M=B\mathrm{\Psi }_M^{},`$ with $`BB^{}=1`$. However, if $`BB^{}=1`$ one can define $`SU(2)`$Majorana spinors (or $`SU(2)`$pseudo-Majorana spinors). This is a pair of spinors $`\mathrm{\Psi }_i,i=1,2`$ satisfying
$`(\mathrm{\Psi }_i)^{}=ϵ^{ij}B\mathrm{\Psi }_j,`$ (4.40)
where $`ϵ^{ij}`$ is the $`SU(2)`$invariant antisymmetric tensor, and $`i,j=1,2`$. (More generally one can take an even number of spinors and substitute to $`ϵ`$ the symplectic form $`\mathrm{\Omega }`$. These spinors are also called symplectic spinors.) The $`SU(2)`$Majorana spinors exist when $`ϵ_1=1,\eta _1=1`$ or $`ϵ_2=1,\eta _2=1`$ and the the $`SU(2)`$pseudo-Majorana spinors exist when $`ϵ_1=1,\eta _1=1`$ or $`ϵ_2=1,\eta _2=1`$
###### Proposition 4.4
Assume $`d`$ even.
1. $`SU(2)`$Majorana spinors of $`\text{P}in(1,d1)`$ exist when $`d=4,6\text{ mod. }8`$.
2. $`SU(2)`$pseudo-Majorana spinors of $`\text{P}in(1,d1)`$ exist when $`d=6,8\text{ mod. }8`$.
Majorana-Weyl and $`SU(2)`$Majorana-Weyl spinors, $`d`$ even
Now, from (4.19) we observe that the Weyl condition is compatible with the Majorana or the $`SU(2)`$Majorana condition if $`n+1=2,4\text{ mod. }4`$. Indeed, in such space-time dimension $`\chi =(1)^{\left[\frac{n+1}{2}\right]}\mathrm{\Gamma }_0\mathrm{}\mathrm{\Gamma }_{2n1}`$ and the chirality matrix is real. This means in this case that the Weyl spinors $`\mathrm{\Psi }_L`$ and $`\mathrm{\Psi }_R`$ can be taken real or $`SU(2)`$Majorana. This is possible only when the space-time dimension $`d=2,6\text{ mod. }8`$. Such spinor will be called Majorana-Weyl spinors ($`d=2\text{ mod. }8`$) and $`SU(2)`$Majorana-Weyl ($`d=6\text{ mod. }8`$). In fact this is the dimensions for which $`𝒞_0_{1,d1}`$ is not simple see Proposition 2.5 (iv).
###### Proposition 4.5
Assume $`d`$ even.
* Majorana-Weyl spinors of $`\text{P}in(1,d1)`$ exist when $`d=2\text{ mod. 8}`$.
* $`SU(2)`$Majorana-Weyl spinors of $`\text{P}in(1,d1)`$ exist when $`d=6\text{ mod. 8}`$.
Space-time of odd dimension $`d=2n+1`$
When the space-time is odd, the matrices $`\mathrm{\Gamma }_\mu ,\mu =0,\mathrm{},2n1`$ are the same as the ones for $`d=2n`$ together with
$`\mathrm{\Gamma }_{2n}=i\chi .`$ (4.41)
In this case taking the matrices (4.11) there is no matrix $`U`$ such that $`U\mathrm{\Gamma }_\mu U^1=\mathrm{\Gamma }_\mu `$, because when $`d`$ is odd Clifford algebras are not simple and $`\gamma (ϵ)`$ is a number. This means in particular that differently to the case where $`d`$ is even we will have either $`B\mathrm{\Gamma }_\mu B^1=\mathrm{\Gamma }_\mu ^{}`$ or $`B\mathrm{\Gamma }_\mu B^1=\mathrm{\Gamma }_\mu ^{}`$ where the sign depends on the dimension (similar property holds for the $`C`$ matrix). The analogous of (4.32) writes:
$`\begin{array}{cc}A\mathrm{\Gamma }_\mu A^1=\mathrm{\Gamma }_\mu ^{},\hfill & A=\mathrm{\Gamma }_0\hfill \\ B\mathrm{\Gamma }_\mu B^1=(1)^n\mathrm{\Gamma }_\mu ^{},\hfill & B=\mathrm{\Gamma }_2\mathrm{\Gamma }_4\mathrm{}\mathrm{\Gamma }_{2n}\hfill \\ C\mathrm{\Gamma }_\mu C^1=(1)^n\mathrm{\Gamma }_\mu ^t,\hfill & C=\mathrm{\Gamma }_0\mathrm{\Gamma }_2\mathrm{}\mathrm{\Gamma }_{2n}\hfill \end{array}`$ (4.45)
with
$`BB^{}=(1)^{\frac{n(n1)}{2}}.`$ (4.46)
One can check that the matrices $`B^{}=\mathrm{\Gamma }_0\mathrm{\Gamma }_1\mathrm{}\mathrm{\Gamma }_{2n1}`$ and $`C^{}=\mathrm{\Gamma }_1\mathrm{\Gamma }_3\mathrm{}\mathrm{\Gamma }_{2n1}`$ give the same relations (4.45) as the matrices $`B,C`$. As for the even space-time dimension, we introduce $`ϵ=(1)^{\frac{n(n1)}{2}},\eta =(1)^n`$ which gives
$`\begin{array}{ccc}d=1\text{ mod. }8\hfill & ϵ=+\hfill & \eta =+\hfill \\ d=3\text{ mod. }8\hfill & ϵ=+\hfill & \eta =\hfill \\ d=5\text{ mod. }8\hfill & ϵ=\hfill & \eta =+\hfill \\ d=7\text{ mod. }8\hfill & ϵ=\hfill & \eta =.\hfill \end{array}`$ (4.51)
Then as for even dimensional space-time we have the following:
###### Proposition 4.6
Assume $`d`$ is odd.
1. Majorana spinors of $`\text{Pin}(1,d1)`$ exist when $`d=1\text{ mod. }8`$.
2. Pseudo-Majorana spinors of $`\text{Pin}(1,d1)`$ exist when $`d=3\text{ mod. }8`$
3. $`SU(2)`$ Majorana spinors of $`\text{Pin}(1,d1)`$ exist when $`d=5\text{ mod. }8`$.
4. $`SU(2)`$ pseudo-Majorana spinors of $`\text{Pin}(1,d1)`$ exist when $`d=7\text{ mod. }8`$.
We conclude this subsection by the following table.
### 4.3 Real, pseudo-real and complex representations of $`𝔰𝔬(1,d1)`$
Now a group theory touch can be given to the different type of spinors. As we have seen, we have the following inclusion of algebras:
$`𝔰𝔬(1,d1)𝒞_{1,d1}.`$ (4.52)
Recall that the generators of $`𝔰𝔬(1,d1)`$ are given by $`\mathrm{\Gamma }_{\mu \nu }=\frac{1}{4}(\mathrm{\Gamma }_\mu \mathrm{\Gamma }_\nu \mathrm{\Gamma }_\nu \mathrm{\Gamma }_\mu ),\mu ,\nu =0,\mathrm{},d1`$,
and the representation $`\mathrm{\Psi }_D`$ and $`B^1\mathrm{\Psi }_D^{}`$ are equivalent. As we have seen they can be equated if $`BB^{}=1`$. The representation of $`𝔰𝔬(1,d1)`$ is called real if they can be equated and pseudo-real when they cannot. In odd dimensional space-time the representations can be either real or pseudo-real. In even dimension in addition to (pseudo-)real representations there exist complex representations.
* When $`d=2n+1`$, the representation of $`𝔰𝔬(1,2n)`$ are real when $`d=1,3\text{ mod. }8`$ and pseudo-real when $`d=5,7\text{ mod. }8`$.
* When $`d=2n`$, we introduce
$`\mathrm{\Sigma }_{\mu \nu }^\pm ={\displaystyle \frac{1}{2}}(1\pm \chi )\mathrm{\Gamma }_{\mu \nu }`$ (4.53)
which are the generators of the spinor representations $`𝒮_{}`$. From (4.19) we get $`\chi ^{}=(1)^{n+1}B\chi B^1,`$ thus
$`\mathrm{\Sigma }_{\mu \nu }^\pm {}_{}{}^{}=\{\begin{array}{cc}B\mathrm{\Sigma }_{\mu \nu }^\pm B^1\hfill & \text{when }n\text{ odd}\hfill \\ B\mathrm{\Sigma }_{\mu \nu }^{}B^1\hfill & \text{when }n\text{ even}.\hfill \end{array}`$ (4.56)
This means that the complex conjugate of $`𝒮_\pm `$ is equal to itself when $`d=4k+2`$. The representation will be real when $`d=2\text{ mod. }8`$ ($`BB^{}=1`$) and pseudo-real when $`d=6\text{ mod. }8`$ ($`BB^{}=1`$). However, when $`d=4k`$ the complex conjugate of $`𝒮_\pm `$ is $`𝒮_{}`$ and the representation is complex.
###### Remark 4.7
The results above can also give some insight on the structure of Clifford algebras. When $`d`$ is even, there always exists a matrix $`B`$ such that $`B\mathrm{\Gamma }_\mu B^1=\mathrm{\Gamma }_\mu `$ (see (4.32)). If $`BB^{}=1`$ the Clifford algebra is real $`(d=0,2\text{ mod. }8`$) and if $`BB^{}=1`$ the algebra is quaternionian ($`d=4,6\text{ mod. }8`$). (Notice that for $`d=4`$ the pseudo-Majorana spinors a taken with the oposite choise i.e with $`B`$ s.t. $`B\mathrm{\Gamma }_\mu B^1=\mathrm{\Gamma }_\mu ^{}`$.) When $`d`$ is odd, either $`\epsilon ^2=1`$ ($`\rho (\epsilon )=\pm `$) or $`\epsilon ^2=1`$ ($`\rho (\epsilon )=\pm i`$). When $`\epsilon ^2=1`$ (or when $`B\mathrm{\Gamma }_\mu B^1=\mathrm{\Gamma }_\mu `$) the Clifford algebra is real when $`BB^{}=1`$ ($`d=1\text{ mod. }8`$) and quaternionian when $`BB^{}=1`$ ($`d=5\text{ mod. }8`$) although when $`\epsilon ^2=1`$ (or when $`B\mathrm{\Gamma }_\mu B^1=\mathrm{\Gamma }_\mu `$) the Clifford algebra is complex $`(d=3,7\text{ mod. }8`$). This can be compared with table 2.
### 4.4 Properties of (anti-)symmetry of the $`\mathrm{\Gamma }`$matrices
The matrices
$`\mathrm{\Gamma }_{\mu _1\mathrm{}\mu _{\mathrm{}}}^{(\mathrm{})}={\displaystyle \frac{1}{\mathrm{}!}}{\displaystyle \underset{\sigma \mathrm{\Sigma }_{\mathrm{}}}{}}ϵ(\sigma )\mathrm{\Gamma }_{\mu _{\sigma (1)}\mathrm{}\mu _{\sigma (\mathrm{})}},`$ (4.57)
with $`\mathrm{}=0,\mathrm{},d`$ when $`d`$ is even and with $`\mathrm{}=0\mathrm{},\frac{d1}{2}`$ when $`d`$ is odd constitute a basis of the representation of $`𝒞_{1,d1}`$. (In $`d`$ dimensions we have $`\left(\begin{array}{c}d\\ \mathrm{}\end{array}\right)`$ independant matrices $`\mathrm{\Gamma }_{\mu _1\mathrm{}\mu _{\mathrm{}}}^{\mathrm{}}`$ and $`2^{2n}=\underset{\mathrm{}=0}{\overset{2n}{}}\left(\begin{array}{c}2n\\ \mathrm{}\end{array}\right)`$ when $`d=2n`$ or $`2^{2n}=\underset{\mathrm{}=0}{\overset{n}{}}\left(\begin{array}{c}2n+1\\ \mathrm{}\end{array}\right)`$ when $`d=2n+1`$.) Note the normalisation factor $`\mathrm{\Gamma }_{\mu \nu }=\frac{1}{2}\mathrm{\Gamma }_{\mu \nu }^{(2)}.`$ Now using the matrices $`C_1`$ and $`C_2`$ in (4.32) and the matrix $`C`$ in (4.45), on can show that the matrices $`\mathrm{\Gamma }^{(\mathrm{})}C^1`$ are either fully symmetric or fully antisymmetric. It is just a matter of a simple calculation to check the following:
$`\begin{array}{cc}(\mathrm{\Gamma }_{\mu _1\mathrm{}\mu _{\mathrm{}}}^{(\mathrm{})}C_1^1)^t=(1)^{\frac{1}{2}(\mathrm{}+n)^2+(\mathrm{}n))}\mathrm{\Gamma }_{\mu _1\mathrm{}\mu _{\mathrm{}}}^{(\mathrm{})}C_1^1\hfill & d=2n\hfill \\ (\mathrm{\Gamma }_{\mu _1\mathrm{}\mu _{\mathrm{}}}^{(\mathrm{})}C_2^1)^t=(1)^{\frac{1}{2}(\mathrm{}+n)^2+(n\mathrm{}))}\mathrm{\Gamma }_{\mu _1\mathrm{}\mu _{\mathrm{}}}^{(\mathrm{})}C_2^1\hfill & d=2n\hfill \\ (\mathrm{\Gamma }_{\mu _1\mathrm{}\mu _{\mathrm{}}}^{(\mathrm{})}C^1)^t=(1)^{\frac{1}{2}(\mathrm{}+n)^2+(n\mathrm{}))}\mathrm{\Gamma }_{\mu _1\mathrm{}\mu _{\mathrm{}}}^{(\mathrm{})}C^1\hfill & d=2n+1\hfill \end{array}`$ (4.61)
We now summarize in the following table the symmetry properties of the $`\mathrm{\Gamma }C^1`$matrices.
Now we would like to give some duality properties. We define $`\epsilon ^{\mu _1\mathrm{}\mu _{2n}}`$ the Levi-Civita tensor (equal to the signature of the permutation that brings $`\mu _1\mathrm{}\mu _{2n}`$ to $`0,1,\mathrm{},2n1,\epsilon ^{0\mathrm{}2n1}=1`$). We also introduce $`\mathrm{\Gamma }^\mu :\mathrm{\Gamma }_\mu =\eta _{\mu \nu }\mathrm{\Gamma }^\nu `$ and $`\epsilon _{\mu _1\mathrm{}\mu _{2n}}=\epsilon ^{\nu _1\mathrm{}\nu _{2n}}\eta _{\mu _1\nu _1}\mathrm{}\eta _{\mu _{2n}\nu _{2n}}`$ $`\epsilon _{0\mathrm{}2n1}=1`$) and $`\mathrm{\Gamma }=\mathrm{\Gamma }_0\mathrm{}\mathrm{\Gamma }_{2n1}`$. From,
$`\mathrm{\Gamma }^{(\mathrm{})}{}_{}{}^{\mu _1\mathrm{}\mu _{2n}}\mathrm{\Gamma }`$ $`=`$ $`(1)^n\epsilon ^{\mu _1\mathrm{}\mu _{2n}}`$ (4.62)
$`\mathrm{\Gamma }^{(\mathrm{})}{}_{}{}^{\mu _1\mathrm{}\mu _{\mathrm{}}}\mathrm{\Gamma }`$ $`=`$ $`{\displaystyle \frac{(1)^{\frac{\mathrm{}(\mathrm{}1)}{2}}}{(2n\mathrm{})!}}\epsilon ^{\mu _1\mathrm{}\mu _{2n}}\mathrm{\Gamma }^{(\mathrm{})}_{\mu _{\mathrm{}+1}\mathrm{}\mu _{2n}}`$
we get
$`\begin{array}{cc}{}_{}{}^{}(\mathrm{\Gamma }^{(n)}\pm i\mathrm{\Gamma }^{(n)}\mathrm{\Gamma })=\pm i(1)^{\frac{n(n+1)}{2}+1}(\mathrm{\Gamma }^{(n)}\pm i\mathrm{\Gamma }^{(n)}\mathrm{\Gamma })\hfill & \text{ when }n\text{ is even}\hfill \\ {}_{}{}^{}(\mathrm{\Gamma }^{(n)}\pm \mathrm{\Gamma }^{(n)}\mathrm{\Gamma })=\pm (1)^{\frac{n(n+1)}{2}+1}(\mathrm{\Gamma }^{(n)}\pm \mathrm{\Gamma }^{(n)}\mathrm{\Gamma })\hfill & \text{ when }n\text{ is odd}\hfill \end{array}`$ (4.65)
with the Hodge dual $`({}_{}{}^{}X_{\mu _1\mathrm{}\mu _p}^{}=\frac{1}{(2np)!}\epsilon _{\mu _1\mathrm{}\mu _p\nu _1\mathrm{}\nu _{2np}}X^{\nu _1\mathrm{}\nu _{2np}}`$). Thus this means that the above matrices are (anti-)self-dual.
### 4.5 Product of spinors
Now we give the decomposition of the product of spinor representations. We recall that $`\text{End}(𝒮)𝒮𝒮^{}`$. Furthermore, given $`\mathrm{\Psi }_D`$ and $`\mathrm{\Xi }_D`$ two Dirac spinors, recall
$`S=\overline{\mathrm{\Xi }}_D\mathrm{\Psi }_D=\mathrm{\Xi }_D^{}\mathrm{\Gamma }_0\mathrm{\Psi }_D`$ (4.66)
define invariant scalar products (see (4.17)).
###### Remark 4.8
When the signature is $`(t,s)`$ we take $`\mathrm{\Gamma }_1\mathrm{}\mathrm{\Gamma }_t`$ instead of $`\mathrm{\Gamma }_0`$ to set $`\overline{\mathrm{\Psi }}=\mathrm{\Psi }^{}\mathrm{\Gamma }_1\mathrm{}\mathrm{\Gamma }_t`$. This legitimate (3.11).
Similarly,
$`\overline{\mathrm{\Xi }}\mathrm{\Gamma }_{\mu _1\mathrm{}\mu _{\mathrm{}}}^{(\mathrm{})}\mathrm{\Psi }_D`$ (4.67)
transforms as an $`\mathrm{}^{\text{th}}`$ order antisymmetric tensor. Now, from (4.17) it is easy to see that
$`T_{\mu _1\mathrm{}\mu _{\mathrm{}}}=\mathrm{\Xi }^tC\mathrm{\Gamma }_{\mu _1\mathrm{}\mu _{\mathrm{}}}^{(\mathrm{})}\mathrm{\Psi }_D`$ (4.68)
transforms also as an $`\mathrm{}^{\text{th}}`$ order antisymmetric tensor. Now, since the matrices $`\mathrm{\Gamma }_{\mu _1\mathrm{}\mu _{\mathrm{}}}^{(\mathrm{})}`$ with $`\mathrm{}=0,\mathrm{}d1`$ when $`d`$ is even and $`\mathrm{}=0,\mathrm{}\frac{d1}{2}`$ when $`d`$ is odd constitute a basis of $`_{2^{\left[\frac{d}{2}\right]}}()`$ and since $`\mathrm{\Xi }^tC𝒮^{}`$ (the dual space of $`𝒮`$) we have
$`𝒮𝒮^{}=\{\begin{array}{cc}E\mathrm{\Lambda }^2(E)\mathrm{}\mathrm{\Lambda }^d(E)\hfill & \text{when }d\text{ is even}\hfill \\ E\mathrm{\Lambda }^3(E)\mathrm{}\mathrm{\Lambda }^{\left[\frac{d}{2}\right]}(E)\hfill & \text{when }d\text{ is odd}.\hfill \end{array}`$ (4.71)
This decomposition has to be compared with Remark 2.4.
Now, when $`d=2n`$ is even one can calculate the product of Weyl spinors. Let $`\mathrm{\Psi }_{ϵ_1}𝒮_{ϵ_1},ϵ_1=\pm `$ and $`\mathrm{\Xi }_{ϵ_2}^tC𝒮_{ϵ_2}^{},ϵ_2=\pm `$ be two Weyl spinors. Recall that we have $`\chi \mathrm{\Psi }_{ϵ_1}=ϵ_1\mathrm{\Psi }_{ϵ_1}`$. Now, using $`C\chi C^1=\mathrm{\Gamma }_0^t\mathrm{}\mathrm{\Gamma }_{d1}^t`$. By (4.25) the RHS writes $`(1)^n\chi `$ but rearranging the product the RHS also writes $`(1)^n\chi ^t`$. Thus
$`\chi ^t=\chi ,C\chi =(1)^n\chi C,`$ (4.72)
and we have
$`\mathrm{\Xi }_{ϵ_2}^tC\mathrm{\Gamma }_{\mu _1\mathrm{}\mu _{\mathrm{}}}^{(\mathrm{})}\mathrm{\Psi }_{ϵ_1}`$ $`=`$ $`ϵ_1\mathrm{\Xi }_{ϵ_2}^tC\mathrm{\Gamma }_{\mu _1\mathrm{}\mu _{\mathrm{}}}^{(\mathrm{})}\chi \mathrm{\Psi }_{ϵ_1}=ϵ_1(1)^{\mathrm{}}\mathrm{\Xi }_{ϵ_2}^tC\chi \mathrm{\Gamma }_{\mu _1\mathrm{}\mu _{\mathrm{}}}^{(\mathrm{})}\mathrm{\Psi }_{ϵ_1}`$ (4.73)
$`=`$ $`ϵ_1(1)^{n+\mathrm{}}\mathrm{\Xi }_{ϵ_2}^t\chi C\mathrm{\Gamma }_{\mu _1\mathrm{}\mu _{\mathrm{}}}^{(\mathrm{})}\mathrm{\Psi }_{ϵ_1}=ϵ_1(1)^{n+\mathrm{}}(\chi \mathrm{\Xi }_{ϵ_2})^tC\mathrm{\Gamma }_{\mu _1\mathrm{}\mu _{\mathrm{}}}^{(\mathrm{})}\mathrm{\Psi }_{ϵ_1}`$
$`=`$ $`ϵ_1ϵ_2(1)^{n+\mathrm{}}\mathrm{\Xi }_{ϵ_2}^tC\mathrm{\Gamma }_{\mu _1\mathrm{}\mu _{\mathrm{}}}^{(\mathrm{})}\mathrm{\Psi }_{ϵ_1}`$
and $`\mathrm{\Xi }_{ϵ_2}^tC\mathrm{\Gamma }_{\mu _1\mathrm{}\mu _{\mathrm{}}}^{(\mathrm{})}\mathrm{\Psi }_{ϵ_1}=0`$ if $`ϵ_1ϵ_2(1)^{n+\mathrm{}}=1`$. This finally gives
$`𝒮_+𝒮_+^{}=\{\begin{array}{cc}\mathrm{\Lambda }^2(E)\mathrm{}\mathrm{\Lambda }^n(E)_+\hfill & \text{ when }n\text{ is even}\hfill \\ E\mathrm{\Lambda }(E)\mathrm{}\mathrm{\Lambda }^n(E)_+\hfill & \text{ when }n\text{ is odd}\hfill \end{array}`$ (4.76)
(4.77)
$`𝒮_+𝒮_{}^{}=\{\begin{array}{cc}E\mathrm{\Lambda }^3(E)\mathrm{}\mathrm{\Lambda }^{n1}(E)\hfill & \text{ when }n\text{ is even}\hfill \\ \mathrm{\Lambda }^2(E)\mathrm{}\mathrm{\Lambda }^{n1}(E)\hfill & \text{ when }n\text{ is odd}\hfill \end{array}.`$ (4.80)
A special attention has to be paid for the antisymmetric tensors of order $`n`$. Indeed, the antisymmetric tensor of order $`n`$ is a reducible representation of $`SO(1,2n1)`$ and can be decomposed into self-dual and anti-self-dual tensors (see (4.65)) (denoted $`\mathrm{\Lambda }^n(E)_\pm `$):
$`\mathrm{\Lambda }^n(E)=\mathrm{\Lambda }^n(E)_+\mathrm{\Lambda }^n(E)_{}.`$ (4.81)
Now a simple counting of the dimensions on both sides in (4.76) shows that only the self- (or anti-self-) dual $`n^{\text{th}}`$order tensors appears in the decomposition. (The dimension of $`\mathrm{\Lambda }^p(E)`$ is $`\left(\begin{array}{c}2n\\ p\end{array}\right)`$ and of $`\mathrm{\Lambda }^n(E)_\pm `$ is $`\frac{1}{2}\left(\begin{array}{c}2n\\ n\end{array}\right)`$.) The choice $`\mathrm{\Lambda }^n(E)_+`$ is a matter of convention.
If we rewrite explictly these relation (e.g. (4.71) for $`d`$ even), we have
$`\mathrm{\Psi }_D\mathrm{\Xi }_D^tC`$ $`=`$ $`{\displaystyle \underset{\mathrm{}=0}{\overset{d}{}}}{\displaystyle \frac{1}{\mathrm{}!}}T^{(\mathrm{})}{}_{}{}^{\mu _1\mathrm{}\mu _{\mathrm{}}}\mathrm{\Gamma }_{\mu _1\mathrm{}\mu _{\mathrm{}}}^{(\mathrm{})}`$
$`\mathrm{\Psi }_D\mathrm{\Xi }_D^t`$ $`=`$ $`{\displaystyle \underset{\mathrm{}=0}{\overset{d}{}}}{\displaystyle \frac{1}{\mathrm{}!}}T^{(\mathrm{})}{}_{}{}^{\mu _1\mathrm{}\mu _{\mathrm{}}}\mathrm{\Gamma }_{\mu _1\mathrm{}\mu _{\mathrm{}}}^{(\mathrm{})}C^1`$ (4.82)
with $`T^{(\mathrm{})}`$ an antisymmetric $`\mathrm{}`$th order tensor.
## 5 Clifford algebras and supersymmetry
In this section, with the help of the previous sections, we study supersymmetric algebras with a special attention to the four, ten and eleven-dimensional space-times. We then study the representations of the considered supersymmetric algebras and show that representation spaces contain an equal number of bosons and fermions. Supersymmetry turns out to be a symmetry which mixes non-trivially the bosons and the fermions since one multiplet contains bosons and fermions together. We also show how four and ten dimensional supersymmetry are related to eleven dimensional supersymmetry by compactification or dimensional reduction.
### 5.1 Non-trivial extensions of the Poincaré algebra
Describing the laws of physics in terms of underlying symmetries has always been a powerful tool. For instance the Casimir operators of the Poincaré algebra (3.1) are related to the mass and the spin of elementary particles as the electron or the photon. Moreover, it has been understood that all the fundamental interactions (electromagnetic, week and strong interactions) are related to the Lie algebra $`𝔲(1)_Y\times 𝔰𝔲(2)_L\times 𝔰𝔲(3)_c`$, in the so-called standard model (see e.g. and references therein). The standard model is then described by the Lie algebra $`𝔦𝔰𝔬(1,3)\times 𝔲(1)_Y\times 𝔰𝔲(2)_L\times 𝔰𝔲(3)_c`$, where $`𝔦𝔰𝔬(1,3)`$ is related to space-time symmetries and $`𝔲(1)_Y\times 𝔰𝔲(2)_L\times 𝔰𝔲(3)_c`$ to internal symmetries. Even if the standard model is the physical theory were the confrontation between experimental results and theoretical predictions is in an extremely good accordance, there is strong arguments (which cannot be summarized here) that it is not the final theory.
Thus, to understand the properties of elementary particles, it is then interesting to study the kind of symmetries which are allowed in space-time. Within the framework of Quantum Field Theory (unitarity of the $`S`$ matrix etc.), S. Coleman and J. Mandula have shown that if the symmetries are described in terms of Lie algebras, only trivial extensions of the Poincaré algebra can be obtained. Namely, the fundamental symmetries are based on $`𝔦𝔰𝔬(1,3)\times 𝔤`$ with $`𝔤𝔲(1)_Y\times 𝔰𝔲(2)_L\times 𝔰𝔲(3)_c`$ a compact Lie algebra describing the fundamental interactions and $`[𝔦𝔰𝔬(1,3),𝔤]=0`$. Several algebras, in relation to phenomenology, have been investigated (see e.g. ) such as $`𝔰𝔲(5),𝔰𝔬(10),𝔢_6`$ etc. Such theories are usually refer to “Grand-Unified-Theories” or GUT, i.e. theories which unify all the fundamental interactions. The fact that elements of $`𝔤`$ and $`𝔦𝔰𝔬(1,3)`$ commute means that we have a trivial extension of the Poincaré algebra.
Then, R. Haag, J. T. Lopuszanski and M. F. Sohnius understood that is was possible to extend in a non-trivial way the symmetries of space-time within the framework of Lie superalgebras (see Definition 5.1) in an unique way called supersymmetry. We first give the definition of a Lie superalgebras, and then we show how to construct a supersymmetric theory in any space-time dimensions using the results established in Section 4.
###### Definition 5.1
A Lie (complex or real) superalgebra is a $`_2`$graded vector space $`𝔤=𝔤_0𝔤_1`$ endowed with the following structure
1. $`𝔤_0`$ is a Lie algebra, we denote by $`[,]`$ the bracket on $`𝔤_0`$ ($`[𝔤_0,𝔤_0]𝔤_0`$);
2. $`𝔤_1`$ is a representation of $`𝔤_0`$ ($`[𝔤_0,𝔤_1]𝔤_1`$);
3. there exits a $`𝔤_0`$equivariant mapping $`\{,\}:S^2\left(𝔤_1\right)𝔤_0`$ where $`S^2\left(𝔤_1\right)`$ denotes the two-fold symmetric product of $`𝔤_1`$ ($`\{𝔤_1,𝔤_1\}𝔤_0`$);
4. The following Jacobi identities hold ($`b_1,b_2,b_3𝔤_0,f_1,f_2,f_3𝔤_1`$)
$`[[b_1,b_2],b_3]+[[b_2,b_3],b_1]+[[b_3,b_1],b_2]=0`$
$`[[b_1,b_2],f_3]+[[b_2,f_3],b_1]+[[f_3,b_1],b_2]=0`$
$`[b_1,\{f_2,f_3\}]\{[b_1,f_2],f_3\}\{f_2,[b_1,f_3]\}=0`$
$`[f_1,\{f_2,f_3\}]+[f_2,\{f_3,f_1\}]+[f_3,\{f_1,f_2\}]=0.`$
The generators of zero (resp. one) gradation are called the bosonic (resp. fermionic) generators or $`𝔤_0`$ (resp. $`𝔤_1`$) is called the bosonic (resp. fermionic) part of the Lie superalgebra. The first Jacobi identity is the usual Jacobi identity for Lie algebras, the second says that $`𝔤_1`$ is a representation of $`𝔤_0`$, the third identity is the equivariance of $`\{,\}`$. These identities are just consequences of $`1.,2.`$ and $`3.`$ respectively. However, the fourth Jacobi identity which is an extra constraint, is just the $`_2`$graded Leibniz rule.
The supersymmetric extension of the Poincaré algebra is constructed, in the framework of Lie superalgebras, by adjoining to the Poincaré generators anticommuting elements, called supercharges (we denote $`Q`$), which belong to the spinor representation of the Poincaré algebra. Thus the supersymmetric algebra is a Lie superalgebra $`𝔤=𝔦𝔰𝔬(1,d1)𝒮`$ with brackets
$`[L,L]=L,[L,P]=P,[L,Q]=Q,[P,Q]=0,\{Q,Q\}=P,`$ (5.2)
with $`(L,P)`$ the generators of the Poincaré algebra that belong to the bosonic part of the algebra and $`Q`$ the fermionic part of the algebra. This extension is non-trivial, because the supercharges $`Q`$ are spinors, and thus do not commute with the generators of the Lorentz algebra. However, the precise definition depends on the space-time dimension because the reality properties of the spinor charges and hence the structure of the algebra depends on the dimensions (see table 3). A systematic study of supersymmetric extensions has been undertaken in . Table 3 indicates the type of spinor we take in various dimensions. Furthermore, the number of supercharges $`N`$ we consider is such that the total spinorial degrees of freedom is less than $`32`$ (see section 5.3.1, Remark 5.5), This gives the possible choices for the spinorial charges (see Table 3)
$`\begin{array}{ccc}d=4\hfill & N\text{pseudo-Majorana}\hfill & 1<N8\hfill \\ d=5\hfill & NSU(2)\text{-Majorana}\hfill & 1<N4\hfill \\ d=6\hfill & (N_+,N_{})\text{ (right-, left-) handed }\hfill & \\ & SU(2)\text{-Majorana-Weyl}\hfill & 1<N_++N_{}4\hfill \\ d=7\hfill & NSU(2)\text{-pseudo-Majorana}\hfill & 1<N2\hfill \\ d=8\hfill & N\text{Majorana}\hfill & 1<N2\hfill \\ d=9\hfill & N\text{Majorana}\hfill & 1<N2\hfill \\ d=10\hfill & (N_+,N_{})\text{ (right-, left-) handed }\hfill & \\ & \text{ Majorana-Weyl}\hfill & 1<N_++N_{}2\hfill \\ d=11\hfill & N\text{pseudo-Majorana}\hfill & N=1\hfill \end{array}`$ (5.13)
### 5.2 Algebra of supersymmetry in various dimensions
Using (5.13) we give the supersymmetric algebras in four, ten and eleven dimensions. Supersymmetry in other space-time dimensions are constructed along the same lines . In this section, we will not give precise references of the subject, one may for instances see and references therein (in particular we will not refer to the original papers on the subject<sup>4</sup><sup>4</sup>4Supersymmerty and supergravity is an intense subject of research. If one goes to the particle physics data basis http://www.slac.standford.edu and types find title supersymmetry or sypergravity, one has 8767 different publications or types find k supersymmetry or supergravity, one has 42449 answers (the 1.06.2005) ).
#### 5.2.1 Supersymmetry in four dimensions
Now, we give the precise structure of the supersymmetric extension of the Poincaré algebra in four dimensions. A standard reference on the subject is (see also ). The bosonic (or even) part of the algebra is given by the Poincaré generators $`L_{\mu \nu },P_\mu `$. The fermionic (or odd) sector is constituted of $`N`$ Majorana spinor supercharges $`Q_I,I=1,\mathrm{},N`$ (see table 3). The algebraic structure is given by three types of bracket: (i) $`[\text{even },\text{ even}]`$, (ii) $`[\text{even },\text{ odd}]`$ and (iii) $`\{\text{odd },\text{ odd}\}`$, where even/odd means bosonic/fermionic generators. The first types of bracket is the Poincaré algebra (3.1). The action of the Poincaré algebra onto the fermionic supercharges is given by
$`[L_{\mu \nu },Q_I]=\mathrm{\Gamma }_{\mu \nu }Q_I,[P_\mu ,Q_I]=0,`$ (5.14)
this is the second type of bracket. The first equation is due to the fact that $`Q`$ is in the spinor representation of the Lorentz algebra (with matrix representation $`\mathrm{\Gamma }_{\mu \nu }`$ (4.12)). Since $`P_\mu `$ transforms like a vector, and among the $`\mathrm{\Gamma }`$matrices only $`\mathrm{\Gamma }_\mu `$ transforms as a vector (see (3.20)), the second relation would be $`[P_\mu ,Q]=c\mathrm{\Gamma }_\mu Q`$. But the Jacobi identity (4) involving $`(P,P,Q)`$ leads to $`c=0`$.
Now, we study the last type of brackets involving only odd generators. To be compatible with the literature , we will not use the Dirac $`\mathrm{\Gamma }`$matrices given in Section 4, but
$`\mathrm{\Gamma }_\mu =\left(\begin{array}{cc}0& \sigma _\mu \\ \overline{\sigma }_\mu & 0\end{array}\right)`$ (5.15)
with $`\sigma _\mu `$ the Pauli matrices given in (2.17) and $`(\overline{\sigma }_0,\overline{\sigma }_i)=(\sigma _0,\sigma _i)`$. In this representation the chirality matrix reads $`\chi _4=\left(\begin{array}{cc}\sigma _0& 0\\ 0& \sigma _0\end{array}\right)`$ and the $`B_4`$, $`C_4`$ matrices (4.32) reduce to $`B_4=i\mathrm{\Gamma }_2,`$ $`C_4=i\mathrm{\Gamma }_0\mathrm{\Gamma }_2`$. The Majorana spinor supercharges are defined by
$`Q_I=\left(\begin{array}{c}Q_L_I\\ Q_R_I\end{array}\right)`$ (5.16)
with $`Q_L_I`$ (resp. $`Q_R_I`$) complex left-handed (resp. right-handed) Weyl spinors. The condition $`Q_I^{}=B_4Q_I`$ gives $`Q_{RI}^{}=i\sigma _2Q_L_I`$ and the complex conjugate of a right-handed spinor is a left-handed spinor as we have seen in section 4.3. Finally, remember that in four dimensions we have for the product of two spinors (see (4.5)) $`\mathrm{\Psi }\mathrm{\Xi }^t=T^{(0)}C_4^1+T^{(1)}{}_{}{}^{\mu }\mathrm{\Gamma }_{\mu }^{}C_4^1+\frac{1}{2}T^{(2)}{}_{}{}^{\mu \nu }\mathrm{\Gamma }_{\mu \nu }^{}C_4^1+\frac{1}{6}T^{(3)}{}_{}{}^{\mu \nu \rho }\mathrm{\Gamma }_{\mu \nu \rho }^{}C_4^1+\frac{1}{24}T^{(4)}{}_{}{}^{\mu \nu \rho \sigma }\mathrm{\Gamma }_{\mu \nu \rho \sigma }^{}C_4^1`$.
To construct the last type of brackets, we make the following remarks:
1. the fermionic part of the algebra as to close onto the bosonic part, which reduces to the Lorentz generators $`L_{\mu \nu }`$ and to the generators of the space-time translations $`P_\mu `$;
2. the only symmetric Dirac $`\mathrm{\Gamma }`$ matrices are $`\mathrm{\Gamma }_\mu `$ and $`\mathrm{\Gamma }_{\mu \nu }`$ (see table 4 and equation (4.5));
3. the Jacobi identity (4) as to be satisfied.
The points 1 and 2 give
$$\{Q_I,Q_J^t\}=\delta _{IJ}(aP^\mu \mathrm{\Gamma }_\mu C_4^1+bL^{\mu \nu }\mathrm{\Gamma }_{\mu \nu }C_4^1)$$
with $`Q_J^t`$ the transpose of $`Q_J`$ (in fact instead of $`\delta _{IJ}`$ we could have obtained a symmetric second order tensor, which can always be diagonalised .) Now, the Jacobi identity involving $`(Q,Q,P)`$ gives $`b=0`$ since $`L_{\mu \nu }`$ and $`P_\mu `$ do not commute. Finally for conventional reason, we chose $`a=2`$.
Now we observe that the bosonic part of the algebra can be enlarged by introducing some new (real) scalar generators $`X_{IJ},Y_{IJ}`$ commuting with all the bosonic elements and being antisymmetric $`X_{IJ}=X_{JI},Y_{IJ}=Y_{JI}`$ (the new generators are called central charges). Since, the matrices $`C_4^1`$ and $`i\chi C_4^1`$ are antisymmetric (see table 4) the algebra extends to
$`\{Q_I,Q_J^t\}=2\delta _{IJ}P^\mu \mathrm{\Gamma }_\mu C_4^1+X_{IJ}C_4^1+iY_{IJ}\chi C_4^1.`$ (5.17)
Now, studying the various Jacobi identities involving $`X,Y`$ we can show
$`[X_{IJ},\text{ anything }]=0,[Y_{IJ},\text{ anything }]=0.`$ (5.18)
The Lie superalgebra defined by (3.1), (5.14), (5.17), and (5.18) is called the four-dimensional $`N`$extended super-Poincaré algebra.
###### Remark 5.2
Since, we have $`N`$ copies of the (complex) supercharges $`Q_L_I`$ and $`Q_R_I`$, this allows the action to the automorphism group for which $`Q_I_L`$ are in the $`N`$dimensional representation of $`GU(N)`$ and $`Q_R_I`$ is in the complex conjugate representation. The full $`N`$extended superalgebra as to be supplemented with $`[T_a,Q_L{}_{I}{}^{}]=(t_a)_I{}_{}{}^{J}Q_{L}^{}{}_{J}{}^{},[T_a,Q_R{}_{I}{}^{}]=(t_a)^{}{}_{I}{}^{}{}_{}{}^{J}Q_{R}^{}_J`$ with $`T_a`$ the generators of $`𝔤`$ (the Lie algebra of $`G`$) and $`t_a`$ the $`N\times N`$ matrices corresponding to the $`N`$dimensional representation of $`𝔤`$ and $`(t_a)^{}`$ the matrices of the complex conjugate representation . With $`Q_I=\left(\begin{array}{c}Q_{LI}\\ Q_{RI}\end{array}\right)`$ these relations become $`[T_a,Q_I]=\frac{(t_a){}_{I}{}^{}{}_{}{}^{J}+(t_a^{}){}_{I}{}^{}^J}{2}Q_J+i\chi \frac{i((t_a){}_{I}{}^{}{}_{}{}^{J}(t_a^{}){}_{I}{}^{}{}_{}{}^{J})}{2}Q_J`$.
###### Remark 5.3
The algebra was presented in a formalism where the spinors have four components. This algebra can also be realized in the two components notations, with the Weyl spinors $`Q_L`$ and $`Q_R`$. From
$`\{Q_I,Q_J^t\}=\{\left(\begin{array}{c}Q_L_I\\ Q_R_I\end{array}\right),\left(\begin{array}{cc}Q_{LJ}^t& Q_{RJ}^t\end{array}\right)\}=\left(\begin{array}{cc}\left\{Q_L{}_{I}{}^{},Q_{LJ}^t\right\}& \left\{Q_L{}_{I}{}^{},Q_{RJ}^t\right\}\\ \left\{Q_R{}_{I}{}^{},Q_{LJ}^t\right\}& \left\{Q_R{}_{I}{}^{},Q_{RJ}^t\right\}\end{array}\right).`$ (5.19)
the algebra (5.17) reduces to
$`\left\{Q_L{}_{I}{}^{},Q_{LJ}^t\right\}=i(X_{IJ}+iY_{IJ})\sigma _2,\left\{Q_L{}_{I}{}^{},Q_{RJ}^t\right\}=2i\delta _{IJ}P^\mu \sigma _\mu \sigma _2.`$
These brackets could also have been deduced from (4.76).
The study of representation of the supersymmetric algebra (see sect 5.3.1) will in fact give $`N8`$. This means that the maximum number of fermionic degrees of freedom is $`32(=4\times 8)`$.
#### 5.2.2 Supersymmetry in ten dimensions
There are various ten dimensional supersymmetric algebras, see and references therein. Recall that when $`d=10`$ we can define Majorana-Weyl spinors (see table 3). Such a spinor has $`16`$ components. The structure of the supersymmetric algebra is very similar to the four dimensional case. We just give here, the $`\{\text{odd},\text{odd}\}`$ part of the algebra. Let $`Q`$ be a Majorana spinor in ten dimensions. Then, if we introduce a real central charge $`Z`$ in addition to the Poincaré generators, using Table 4, and arguments similar as in the previous section we get the supersymmetric algebra in dimension ten (since $`\mathrm{\Gamma }_\mu C_{10}^1`$ and $`C_{10}^1`$ are symmetric see table 4)
$`\{Q,Q^t\}=ZC_{10}^1+P^\mu \mathrm{\Gamma }_\mu C_{10}^1,`$ (5.20)
with $`C_{10}`$ the $`C`$matrix in dimension $`10`$ (see (4.32)). From this equation, if we denote $`Q_\pm =\frac{1}{2}(1\pm \chi _{10})Q`$ the left- and right-handed components of $`Q`$, we obtain
$$\{Q_\pm ,Q_\pm ^t\}=\frac{1}{2}(1\pm \chi _{10})P^\mu \mathrm{\Gamma }_\mu C_{10}^1,\{Q_\pm ,Q_{}^t\}=\frac{1}{2}(1\pm \chi _{10})ZC_{10}^1,$$
since $`\chi _{10}C_{10}^1=C_{10}^1\chi _{10}`$ (see eq.(4.72)) and $`\chi _{10}\mathrm{\Gamma }_\mu =\mathrm{\Gamma }_\mu \chi _{10}`$ (see eq.(4.20)) with $`\chi _{10}`$ the chirality (4.19) matrix in ten dimensions. In ten dimensions the fermionic part of the algebra is constituted of $`N_+`$ left-handed Majorana-Weyl spinors and $`N_{}`$ right-handed Majorana-Weyl spinors. As in four dimensions the number of fermionic degrees of freedom is at most $`32`$. Since a Majorana-Weyl spinor has $`16`$ components $`N_++N_{}2`$ leading to three different theories:
1. Type I supersymmetry We have one Majorana-Weyl supercharge, say $`Q_+`$ :
$`\{Q_+,Q_+^t\}={\displaystyle \frac{1}{2}}(1+\chi _{10})P^\mu \mathrm{\Gamma }_\mu C_{10}^1.`$ (5.21)
2. Type IIA supersymmetry We have two Majorana-Weyl supercharges, of opposite chirality $`Q_+`$ and $`Q_{}`$ (or one Majorana spinors $`Q=\left(\begin{array}{c}Q_+\\ Q_{}\end{array}\right)`$) :
$`\{Q_\pm ,Q_\pm ^t\}={\displaystyle \frac{1}{2}}(1\pm \chi _{10})P^\mu \mathrm{\Gamma }_\mu C_{10}^1,\{Q_+,Q_{}^t\}={\displaystyle \frac{1}{2}}(1+\chi _{10})ZC_{10}^1.`$ (5.22)
3. Type IIB supersymmetry We have two Majorana-Weyl supercharges, of the same chirality $`Q_+_1`$ and $`Q_+_2`$:
$`\{Q_{+I},Q_{+J}^t\}=\delta _{IJ}{\displaystyle \frac{1}{2}}(1+\chi _{10})P^\mu \mathrm{\Gamma }_\mu C_{10}^1.`$ (5.23)
The type I (resp. IIA, IIB) theories appears naturally in string theory .
#### 5.2.3 Supersymmetry in eleven dimensions
Now, we give the eleven-dimensional supersymmetric algebra (see e.g. ). From table 3, in eleven dimensions, we take a pseudo-Majorana spinor. Since such a spinor has $`32`$ components there is only one possible theory. Thus, in eleven dimension, the situation is even simpler than in ten or four dimensions. Let $`Q`$ be the pseudo-Majorana supercharge. If the bosonic sector is constituted only of the Poincaré generators, using table 4 the algebra is given by
$`\{Q,Q^t\}=P^\mu \mathrm{\Gamma }_\mu C_{11}^1,`$ (5.24)
with $`C_{11}`$ the eleven-dimensional $`C`$matrix (see (4.45)). The interesting point of the eleven-dimensional supersymmetry is its simplicity. In addition, as we will see below lower dimensional theory can be obtained by dimensional reduction or compactification. Furthermore, eleven is the maximum dimension (with a signature $`(1,d1)`$) where a supersymmetric theory can be formulated. Indeed, when $`d=12`$ the Majorana spinors have $`64`$ components<sup>5</sup><sup>5</sup>5 In twelve dimensions with signature $`(2,10)`$ a Majorana-Weyl spinor has $`32`$ components..
Finally, looking to table 4, one observes that some other types of central charges can be added to extend the algebra (5.24). The possible central charges are real antisymmetric tensors of order two and five leading to the superalgebra
$`\{Q,Q^t\}=P^\mu \mathrm{\Gamma }_\mu C_{11}^1+{\displaystyle \frac{1}{2}}Z_2^{\mu \nu }\mathrm{\Gamma }_{\mu \nu }C_{11}^1+{\displaystyle \frac{1}{5!}}Z_5^{\mu _1\mu _2\mu _3\mu _4\mu _5}\mathrm{\Gamma }_{\mu _1\mu _2\mu _3\mu _4\mu _5}C_{11}^1.`$ (5.25)
The new central charges introduced here are rather different than the central charges considered in four or ten dimensions. Indeed, such central charges are not central, being in antisymmetric tensor representations of the Lorentz algebra they are not scalar. This possibility of tensorial central charges have been considered in . In Quantum Field Theory, or in a theory involving only elementary particles such central charges are excluded by the theorem of Haag-Lopuszanski-Sohnius (since they are not scalars). They become relevant in a theory involving propagating $`p`$dimensional extended objects (strings $`p=1`$, membranes $`p=2`$, etc. called generically branes or $`p`$branes) . For references on $`p`$branes see e.g. . The algebra (5.25) is the basic algebra which underlines $`M`$theory and it involves a membrane and an extended object of dimension five, a $`5`$brane .
Type IIA supersymmetry from eleven-dimensional supersymmetry
If we denote by $`M^{10}`$ and $`M^{11}`$ the Minkowski space-time in $`10`$ or $`11`$ dimensions, we obviously have $`M^{10}M^{11}`$. At the level of the algebra this reduces to $`𝔦𝔰𝔬(1,9)𝔦𝔰𝔬(1,10)`$. This simple observation is the starting point of the so-called dimensional reduction or compactification where a theory in $`10`$ dimensions is obtained from a theory in eleven dimensions. Indeed, historically type IIA supersymmetry was obtained in such a process (see e.g. ). Starting from the eleven dimensional supersymmetry we take the eleventh dimension $`x^{10}`$ to be on a circle $`S^1`$, and we let the radius of the circle to be very small, such that $`M^{11}M^{10}\times S^1`$. If we denote by $`L_{MN},P_M,0M,N10`$ the Poincaré generators in eleven dimensions they reduce to $`L_{\mu \nu },P_\mu ,0\mu ,\nu 9`$ the generators of the $`10`$dimensional Poincaré algebra and to $`P_{10}`$ a scalar with respect to $`𝔦𝔰𝔬(1,9)`$. At the level of Clifford algebras we have $`𝒞_{1,9}𝒞_{1,10}`$. This means that a Majorana spinor in eleven dimensions reduces to two Majorana-Weyl spinors $`Q_+`$ and $`Q_{}`$ in ten dimensions because the Dirac $`\mathrm{\Gamma }`$matrices $`\mathrm{\Gamma }_{10}`$ (presents in eleven dimensions) plays the role of the chirality matrix $`\chi _{10}`$ in ten dimensions (see eq (4.41)). Thus the Majorana spinor reduces to two opposite chirality Majorana-Weyl spinors $`Q_\pm =\frac{1}{2}(1\pm \chi )Q`$. Next observing that the $`C`$matrices in eleven and ten dimensions are related as follow: $`C_{11}=C_{10}\mathrm{\Gamma }_{10}`$, the supersymmetric algebra in ten dimensions obtained from the dimensional reduction of the eleven-dimensional algebra becomes
$$\{Q,Q^t\}=P^\mu \mathrm{\Gamma }_\mu \mathrm{\Gamma }_{10}^1C_{10}^1+P^{10}C_{10}^1.$$
which is just (5.20) with $`Z=P^{10}`$ and $`\mathrm{\Gamma }_\mu \mathrm{\Gamma }_\mu \mathrm{\Gamma }_{10}^1`$. With the two Majorana-Weyl spinors of opposite chirality, we get the type IIA theory in ten dimensions, and $`P^{10}`$ becomes a central charge. This means that the dimensional reduction of the eleven-dimensional supersymmetry leads to the ten dimensional supersymmetry of type IIA.
Four dimensional supersymmetry from eleven-dimensional supersymmetry
The same principle can be applied to construct four dimensional supersymmetry (see for references) from eleven dimensional supersymmetry. But here in the compactification process we have several possibilities for the compact manifold. One of the simplest compactification is to consider a $`7`$sphere such that $`M^{11}M^4\times S^7`$ with $`M^4`$ the four dimensional Minkowski space-time.
In the reduction process from eleven dimensions to four dimensions we observe several things
1. $`𝔰𝔬(1,10)𝔰𝔬(1,3)\times 𝔰𝔬(7)`$;
2. the Poincaré generators $`P^M,M=1,\mathrm{},10`$ give the Poincaré generators in four dimensions $`P^\mu ,\mu =0,\mathrm{},3`$ and $`7`$ scalars with respect to $`𝔦𝔰𝔬(1,3)`$ $`P^4,\mathrm{},P^{10}`$;
3. the spin representation $`\mathrm{𝟑𝟐}`$ of $`\text{ Spin}(1,10)`$ decomposes to $`\mathrm{𝟑𝟐}=(\mathrm{𝟐}_+,\mathrm{𝟖})(\mathrm{𝟐}_{},\mathrm{𝟖})`$ with $`\mathrm{𝟐}_\pm `$ a left/right-handed spinors of $`𝔰𝔬(1,3)`$ and $`\mathrm{𝟖}`$ a Majorana spinor of $`𝔰𝔬(7)`$. (Real spinors exist for $`\text{ Spin}(1,3),\text{ Spin}(7),\text{Spin}(1,10),`$ see table 3.) Thus, the supercharge $`Q`$ gives $`8`$ Majorana supercharges $`Q_I`$ in four dimensions;
4. the Dirac $`\mathrm{\Gamma }`$matrices decompose as follow
$$\mathrm{\Gamma }^M,M=0,\mathrm{},10,\{\begin{array}{cc}\mathrm{\Gamma }_\mu =\mathrm{\Gamma }_\mu ^{(4)}I^{(7)},\hfill & \mu =0,\mathrm{},3,\hfill \\ \mathrm{\Gamma }_m=i\chi _4\mathrm{\Gamma }_m^{(7)},\hfill & m=1,\mathrm{},m\hfill \end{array},$$
with $`\mathrm{\Gamma }_\mu ^{(4)}`$ the four dimensional Dirac $`\mathrm{\Gamma }`$matrices, $`\mathrm{\Gamma }_m^{(7)}`$ the matrices of the representation of $`𝒞_{7,0}`$, $`\chi _4`$ the four dimensional chirality matrix and $`I^{(7)}`$ the identity of $`𝒞_{7,0}`$;
5. the matrix $`C_{11}`$ decomposes into $`C_{11}=C_4C_7`$ with $`C_4=i\mathrm{\Gamma }_0^{(4)}\mathrm{\Gamma }_2^{(4)}`$ and $`C_7=i\mathrm{\Gamma }_4^{(7)}\mathrm{\Gamma }_6^{(7)}\mathrm{\Gamma }_8^{(7)}\mathrm{\Gamma }_{10}^{(7)}`$.
With all these observations the algebra reduces to
$$\{Q,Q^t\}=P^\mu \mathrm{\Gamma }_\mu ^{(4)}C_4^1C_7^1+iP^m\chi _4C_4^1\mathrm{\Gamma }_m^{(7)}C_7^1$$
Next, observing that $`C_7`$ is symmetric (see (4.32)) and introducing $`Q_L{}_{I}{}^{}=(\mathrm{𝟐}_+,\mathrm{𝟖}),Q_R{}_{I}{}^{}=(\mathrm{𝟐}_{},\mathrm{𝟖})`$ and $`Q_I=\left(\begin{array}{c}Q_L_I\\ Q_R_I\end{array}\right)`$ with $`I=1,\mathrm{}8`$ the algebra can be rewritten (after an appropriate diagonalisation $`C_7I^{(7)}`$)
$$\{Q_I,Q_J^t\}=\delta _{IJ}P^\mu \mathrm{\Gamma }_\mu ^{(4)}C_4^1+iP^m\mathrm{\Gamma }_m^{(7)}{}_{I}{}^{}{}_{}{}^{K}\delta _{KJ}^{}\chi _4C_4^1,$$
(with $`\mathrm{\Gamma }_m^{(7)}{}_{I}{}^{}^K`$ the matrix elements of the matrices $`\mathrm{\Gamma }_m^{(7)})`$ which is the $`8`$extended supersymmetric algebra in four dimensions with central charge $`Y_{IJ}=P^m\mathrm{\Gamma }_m^{(7)}{}_{I}{}^{}{}_{}{}^{K}\delta _{KJ}^{}`$ (see (5.17)).
Since the isometry group of $`S^7`$ is $`SO(8)SO(7)`$, in fact in can be shown that the supercharges belong to the vector representation of $`SO(8)`$ (see and references therein).
### 5.3 Irreducible representations of supersymmetry
Since the Poincaré algebra admits a semi-direct structure $`𝔦𝔰𝔬(1,3)=𝔰𝔬(1,3)𝒟_{\frac{1}{2},\frac{1}{2}}`$, with $`𝒟_{\frac{1}{2},\frac{1}{2}}`$ the vector representation of $`𝔰𝔬(1,3)`$ (i.e. the space-time translations), the representations of $`𝔦𝔰𝔬(1,3)`$ are obtained by the method of induced representations of Wigner (this also hold for any space-time dimensions). This method consists of finding a representation of a subgroup of the Poincaré group and boosting it to the full group. If we denote $`P^\mu `$ the momentum, we have $`P_\mu P^\mu =m^2`$ for a particle of mass $`m0`$ and $`P_\mu P^\mu =0`$ for a massless particle<sup>6</sup><sup>6</sup>6 $`P_\mu P^\mu `$ is a Casimir operator of the Poincaré algebra, its eigenvalue is the mass.. In practice, we go in some special frame where the momentum have a specific expression $`P^\mu =(m,0,\mathrm{},0)`$ in the massive case and $`P^\mu =(E,0,\mathrm{},0,E)`$ (with $`E>0`$) in the massless case. We identify the subgroup $`HISO(1,d1)`$, called the little group, which leaves the momentum invariant, find the representations of $`H`$ and then induce the representations to the whole group $`ISO(1,d1)`$. This method, is in fact very close to the principle of equivalence of special relativity which states that if a result is obtained in a specific frame it can be extended to any frame of reference.
In this lecture, we will only study the case of massless particles. If we denote $`L_{\mu \nu }`$ the generators of the Lorentz algebra $`𝔰𝔬(1,d1)`$ the little group leaving $`P^\mu =(E,0,\mathrm{},0,E)`$ $`(P_\mu =P^\nu \eta _{\mu \nu }=(E,0,\mathrm{},0,E)`$) invariant is generated by $`L_{ij},1i<jd2`$ and $`T_i=L_{i0}+L_{id1},1id2`$ (since $`[L_{ij},P_\mu ]=0,[T_i,P_\mu ]=0`$). This group is isomorphic to $`E(d2)`$ the group of rotations-translations in $`(d2)`$ dimensions. Since this group is non-compact and we are interested in finite dimensional unitary representations, we will represent the generators $`T_i`$ by zero. By abuse of notations we will now call $`SO(d2)`$ the little group. The frame where $`P^\mu =(E,0\mathrm{},0,E)`$ will be called the “standard frame”. This method can be extended to study the massless representations of the supersymmetric extension of the Poincaré algebra. The little algebra now contains the bosonic generators $`L_{ij},1i<jd2`$ the central charges (if there exists) and the fermionic supercharges $`Q`$.
#### 5.3.1 Four dimensional supersymmetry
In four dimensions, the generators of the supersymmetric algebra in the little group reduce $`L_{12}`$, $`X_{IJ},Y_{IJ}`$ and $`T_a`$ for the bosonic sector (we assume here that the automorphism group of the algebra is $`SU(N)`$, (see Remark 5.2) and $`Q_I,I=1,\mathrm{},N`$ for the fermionic sector. The generator $`L_{12}`$ is called the generator of helicity. We study here the representations of the supersymmetric algebra when the central charge are put to zero. (In fact it can be proven that massless particles represent trivially the central charges .)
The $`\{\text{odd},\text{odd}\}`$ part of the algebra gives (in the two components notation (see Remark 5.3))
$`\left\{Q_L{}_{I}{}^{},Q_{RJ}^t\right\}=2E\delta _{IJ}(\sigma _0+\sigma _3)i\sigma _2=4E\delta _{IJ}\left(\begin{array}{cc}0& 1\\ 0& 0\end{array}\right).`$ (5.26)
If we write $`Q_L{}_{I}{}^{}=\left(\begin{array}{c}Q_1_I\\ Q_2_I\end{array}\right)`$ and $`Q_R{}_{I}{}^{}=\left(\begin{array}{c}\overline{Q}^{\dot{1}}_I\\ \overline{Q}^{\dot{2}}I\end{array}\right)`$ in the usual notations (see ), the only non-zero brackets are given by
$$\{Q_1{}_{I}{}^{},\overline{Q}^{\dot{2}}{}_{J}{}^{}\}=4E\delta _{IJ}.$$
Since the other brackets vanish and since we want unitary representations this means that $`Q_2{}_{I}{}^{}=\overline{Q}^{\dot{1}}{}_{I}{}^{}=0`$ and the supercharges $`a_I=\frac{Q_1_I}{\sqrt{4E}},a^{}{}_{I}{}^{}=\frac{\overline{Q}^{\dot{2}}_I}{\sqrt{4E}}`$ generate the Clifford algebra $`𝒞_{2N,0}`$ (see Remark 4.1). Now, the action of the bosonic part on the supercharges gives $`[L_{12},Q_I]=\mathrm{\Gamma }_{12}Q_I`$. Using $`\mathrm{\Gamma }_{12}=\frac{1}{2}\mathrm{\Gamma }_1\mathrm{\Gamma }_2`$ and (5.15), we get
$`[L_{12},Q_1{}_{I}{}^{}]=i{\displaystyle \frac{1}{2}}Q_1{}_{I}{}^{},[L_{12},\overline{Q}^{\dot{2}}{}_{I}{}^{}]=i{\displaystyle \frac{1}{2}}\overline{Q}^{\dot{2}}{}_{I}{}^{},`$ (5.27)
and thus $`Q_1_I`$ are of helicity $`\frac{1}{2}`$ and $`\overline{Q}{}_{}{}^{\dot{2}}_I`$ of helicity $`\frac{1}{2}`$. If $`M_{12}|\lambda =i\lambda |\lambda `$ then $`M_{12}Q_1{}_{I}{}^{}|\lambda =i(\lambda +\frac{1}{2})Q_1{}_{I}{}^{}|\lambda `$ and $`M_{12}Q^{\dot{2}}{}_{I}{}^{}|\lambda =i(\lambda \frac{1}{2})Q{}_{}{}^{\dot{2}}{}_{I}{}^{}|\lambda `$. Finally, $`Q_1_I`$ is in the $`N`$dimensional representation of $`SU(N)`$ that we denote $`𝐍`$ and $`\overline{Q}{}_{}{}^{\dot{2}}_I`$ is in the complex conjugate representation $`\overline{𝐍}`$. Thus we have $`Q_1{}_{I}{}^{}=(\frac{1}{2},𝐍),\overline{Q}^{\dot{2}}{}_{I}{}^{}=(\frac{1}{2},\overline{𝐍})`$ with respect to the group $`\text{Spin}(2)\times SU(N)`$.
The representations of the four dimensional supersymmetric algebra are then completely specified and is of dimensions $`2^N`$, corresponding to the spinor representations of $`\text{Spin}(2N)`$. The left-handed spinors of $`\text{Spin}(2N)`$ will correspond e.g. to the fermions and the right-handed spinors to the bosons as we will see. We also know that $`(Q_1{}_{I}{}^{},\overline{Q}^{\dot{2}}{}_{I}{}^{})`$ belongs to the vector representation of $`\text{Spin}(2N)`$. But, if one wants to identify the particles content of the supersymmetric multiplet, it is interesting to decompose the multiplet with respect to the group $`\text{Spin}(2N)\text{Spin}(2)\times SU(N)`$ (i.e. the little group). For the supercharge we have in this embedding $`\mathrm{𝟐}𝐍=(\frac{1}{2},𝐍)(\frac{1}{2},\overline{𝐍})`$ (the supercharge are in the vector representation $`\mathrm{𝟐}𝐍`$ of $`\text{Spin}(2N)`$ -see Remark 4.1-). And using Remark 4.1 it is easy to decompose the spinor representation $`\mathrm{𝟐}^𝐍`$ of $`\text{Spin}(2N)`$ into representations of $`\text{Spin}(2)\times SU(N)`$. If we introduce a Clifford vacuum $`\mathrm{\Omega }`$ of helicity $`\lambda _{\text{max}}`$ ($`L_{12}\mathrm{\Omega }=i\lambda _{\text{max }}\mathrm{\Omega }`$) annihilated by $`a_I(a_I\mathrm{\Omega }=0`$) and being in some representation of $`SU(N)`$ the full representation is obtained by the action of the operator of creation $`a^{}_I`$. For simplification we assume that $`\mathrm{\Omega }`$ is in the trivial representation of $`SU(N)`$ and we obtain the supermultiplet
$`\begin{array}{ccc}\text{state}\hfill & \text{helicity}\hfill & \text{representation of }SU(N)\text{ dimension}\hfill \\ & & \\ |\lambda _{\text{max}}\hfill & \lambda _{\text{max}}\hfill & [0]\text{ dim }=1\hfill \\ \overline{Q}^{\dot{2}}{}_{I}{}^{}|\lambda _{\text{max}}\hfill & \lambda _{\text{max}}\frac{1}{2}\hfill & [\overline{1}]\text{ dim }=N\hfill \\ \overline{Q}^{\dot{2}}{}_{I_1}{}^{}\overline{Q}_{}^{\dot{2}}{}_{I_2}{}^{}|\lambda _{\text{max}}\hfill & \lambda _{\text{max}}1\hfill & [\overline{2}]\text{ dim }=\left(\begin{array}{c}N\\ 2\end{array}\right)\hfill \\ \mathrm{}\hfill & \mathrm{}\hfill & \mathrm{}\hfill \\ \overline{Q}^{\dot{2}}{}_{I_1}{}^{}\overline{Q}_{}^{\dot{2}}{}_{I_2}{}^{}\mathrm{}\overline{Q}^{\dot{2}}{}_{I_k}{}^{}|\lambda _{\text{max}}\hfill & \lambda _{\text{max}}\frac{k}{2}\hfill & [\overline{k}]\text{ dim }=\left(\begin{array}{c}N\\ k\end{array}\right)\hfill \\ \mathrm{}\hfill & \mathrm{}\hfill & \mathrm{}\hfill \\ \overline{Q}^{\dot{2}}{}_{1}{}^{}\overline{Q}_{}^{\dot{2}}{}_{2}{}^{}\mathrm{}\overline{Q}^{\dot{2}}{}_{N}{}^{}|\lambda _{\text{max}}\hfill & \lambda _{\text{max}}\frac{N}{2}\hfill & [\overline{N}]\text{ dim }=1\hfill \end{array}`$ (5.37)
with $`1I_1<I_2<\mathrm{}<I_{N1}N`$, $`[\overline{k}]`$ the antisymmetric tensor of order $`k`$ of $`SU(N)`$ and $`\left(\begin{array}{c}N\\ k\end{array}\right)`$ its dimension. In this decomposition we have
$`(\mathrm{𝟐}^{𝐍\mathrm{𝟏}})_L`$ $`=`$ $`(\lambda _{\text{max }},[0])(\lambda _{\text{max }}1,[2])(\lambda _{\text{max }}2,[4])\mathrm{}`$ (5.38)
$`(\mathrm{𝟐}^{𝐍\mathrm{𝟏}})_R`$ $`=`$ $`(\lambda _{\text{max }}{\displaystyle \frac{1}{2}},[1])(\lambda _{\text{max }}{\displaystyle \frac{3}{2}},[3])\mathrm{}`$
for the left/right-handed part of the spinor representation of $`\text{Spin}(2N)`$. For instance if $`\lambda _{\text{max}}+\frac{1}{2}`$ the left-handed spinors of $`\text{Spin}(2N)`$ are fermions and the right-handed spinors are bosons. A multiplet contains bosons and fermions and a supersymmetric transformation sends a boson to a fermion and vice versa. We also see that we have an equal number of bosons and fermions.
However, these irreducible representations are not enough for particle physics. We still have to take the CPT symmetry into account (C means charge –or complex– conjugation, P parity transformation and T time reversal). Under the CPT symmetry $`Q_RQ_R^{}=i\sigma _2Q_L`$ (or $`a^{}{}_{I}{}^{}a_I`$) and $`L_{12}L_{12}`$. A quantum field theory has to be CPT invariant. This means if the multiplet (5.37) is not invariant under this conjugation, we have to consider the CPT conjugate multiplet obtained by acting on the conjugated Clifford vacuum $`|\lambda _{\text{max}}`$ with the annihilation operator $`Q_1_I`$ ($`\overline{Q}^{\dot{2}}{}_{I}{}^{}|\lambda _{\text{max}}=0`$). Let us give the result for certain values of $`N`$.
1. For $`N=1`$ the multiplet (5.37) are not CPT conjugate. The particle content (after adding the CPT conjugate multiplet) is
1. $`\lambda _{\text{max}}=\frac{1}{2}`$ $`\begin{array}{cccc}\text{helicity}\hfill & \frac{1}{2}\hfill & 0\hfill & \frac{1}{2}\hfill \\ \text{states}\hfill & 1\hfill & 2\hfill & \mathrm{\hspace{0.17em}1}\hfill \end{array}`$;
(In this multiplet we have $`\mathrm{\Omega }=|\frac{1}{2},\overline{Q}^{\dot{2}}\mathrm{\Omega }=|0`$, $`\mathrm{\Omega }^{\text{CPT}}=|\frac{1}{2},Q_1\mathrm{\Omega }^{\text{CPT}}=|0^{}`$.)
2. $`\lambda _{\text{max}}=1`$ $`\begin{array}{ccccc}\text{helicity}\hfill & 1\hfill & \frac{1}{2}\hfill & \frac{1}{2}\hfill & 1\hfill \\ \text{states}\hfill & 1\hfill & 1\hfill & \mathrm{\hspace{0.17em}1}\hfill & \mathrm{\hspace{0.17em}1}\hfill \end{array}`$;
3. $`\lambda _{\text{max}}=2`$ $`\begin{array}{ccccc}\text{helicity}\hfill & 2\hfill & \frac{3}{2}\hfill & \frac{3}{2}\hfill & 2\hfill \\ \text{states}\hfill & 1\hfill & 1\hfill & \mathrm{\hspace{0.17em}1}\hfill & \mathrm{\hspace{0.17em}1}\hfill \end{array}`$.
To identify the particle content of the various multiplet, we have to keep in mind that a massless particle of spin $`s`$ is constituted of a state of helicity $`s`$ and a state of helicity $`s`$. For instance a left-handed massless electron is constituted of a state of helicity $`1/2`$ and a state of helicity $`1/2`$. The former can be interpreted as a left-handed electron of helicity $`\frac{1}{2}`$ and the latter as its corresponding anti-particle state the left-handed positron of helicity $`\frac{1}{2}`$. Thus, after boosting the representations to the whole Poincaré group, the first multiplet contains a left-handed fermion and a complex scalar field and the second multiplet contains a real pseudo-Majorana spinor and a real vector field. These types of multiplet are essential in the construction of the so-called minimal supersymmetric standard model $``$MSSM$``$ (i.e. to construct a model describing particle physics and being supersymmetric). The former multiplets, called the chiral multiplets, are the matter multiplets (they correspond for instance to a left-handed electron and a scalar electron named selectron), although the second multiplets, called the vector multiplets, are relevant for the description of supersymmetric fundamental interactions (in the case of electromagnetism, it corresponds to the photon and its fermionic supersymmetric partner the photino) . This can be generalised for all the particles. This means that in supersymmetric theory the spectrum is doubled, and to each known particle we have to add its supersymmetric partner . The last types of multiplet (the gravity multiplet) contains the graviton (so describes gravity) and its supersymmetric partner, a spinor-vector named the gravitino. It is possible to couple the gravity multiplet with the multiplet of the MSSM .
2. For $`N=4`$, when $`\lambda _{\text{max}}=1`$ the multiplet is CPT conjugate and is not CPT conjugate for $`\lambda _{\text{max}}=2`$:
1. $`\lambda _{\text{max}}=1`$ $`\begin{array}{cccccc}\text{helicity}\hfill & 1\hfill & \frac{1}{2}\hfill & 0\hfill & \frac{1}{2}\hfill & 1\hfill \\ \text{states}\hfill & 1\hfill & 4\hfill & 6\hfill & \mathrm{\hspace{0.17em}4}\hfill & \mathrm{\hspace{0.17em}1}\hfill \end{array}`$;
2.
3. $`\lambda _{\text{max}}=2`$ $`\begin{array}{cccccccccc}\text{helicity}\hfill & 2\hfill & \frac{3}{2}\hfill & 1\hfill & \frac{1}{2}\hfill & 0\hfill & \frac{1}{2}\hfill & 1\hfill & \frac{3}{2}\hfill & 2\hfill \\ \text{states}\hfill & 1\hfill & 4\hfill & 6\hfill & 4\hfill & 2\hfill & \mathrm{\hspace{0.17em}4}\hfill & \mathrm{\hspace{0.17em}6}\hfill & \mathrm{\hspace{0.17em}4}\hfill & \mathrm{\hspace{0.17em}1}\hfill \end{array}`$.
We observe that there is no multiplet with $`\lambda _{\text{max}}=\frac{1}{2}`$ when $`N=4`$. Thus $`N=4`$ does not contain matter multiplets, but only gauge ($`\lambda _{\text{max}}=1`$) or gravity ($`\lambda _{\text{max}}=2`$) multiplets.
3. For $`N=8`$ there is only one (gravity, $`\lambda _{\text{max}}=2`$) multiplet which is CPT conjugate:
1. $`\begin{array}{cccccccccc}\text{helicity}\hfill & 2\hfill & \frac{3}{2}\hfill & 1\hfill & \frac{1}{2}\hfill & 0\hfill & \frac{1}{2}\hfill & 1\hfill & \frac{3}{2}\hfill & 2\hfill \\ \text{states}\hfill & 1\hfill & 8\hfill & 28\hfill & 56\hfill & 70\hfill & 56\hfill & 28\hfill & \mathrm{\hspace{0.17em}8}\hfill & \mathrm{\hspace{0.17em}1}\hfill \end{array}`$
Several remarks are in order here
###### Remark 5.4
One can observe by a direct counting that in a supersymmetric multiplet there is an equal number of bosons and fermions. This is a general result valid in any space-time dimensions (see also (5.38)).
###### Remark 5.5
If $`N>8`$ the supersymmetric multiplets contain states of helicity bigger than $`2`$. Since there is not consistent theory for interacting particles of helicity bigger than $`2`$ in four dimensions, $`N8`$. The complicated multiplets for $`N>1`$ can be obtained by compactification of higher dimensional theories. For instance the four dimensional $`N=8`$ extended supersymmetry can be obtained from the eleven-dimensional theory or the ten dimensional type IIA or type IIB theories and the four dimensional $`N=4`$ extended supersymmetry can be obtained from the ten-dimensional type I theories . Conversely, this compactification limits the number of supersymmetry one can take in a given dimension (see (5.13))
#### 5.3.2 Eleven dimensional supersymmetry
In eleven dimensions, the bosonic part of the little group is $`SO(9)`$ and the supersymmetric algebra becomes
$$\{Q,Q^t\}=E(\mathrm{\Gamma }_0+\mathrm{\Gamma }_{10})C_{11}^1=E(1+\mathrm{\Gamma }_{10}\mathrm{\Gamma }_0)\mathrm{\Gamma }_0C_{11}^1.$$
Since the trace of $`\mathrm{\Gamma }_{10}\mathrm{\Gamma }_0`$ is equal to zero, $`(\mathrm{\Gamma }_{10}\mathrm{\Gamma }_0)^{}=\mathrm{\Gamma }_{10}\mathrm{\Gamma }_0`$ and $`(\mathrm{\Gamma }_{10}\mathrm{\Gamma }_0)^2=1`$, this means that the matrix $`\mathrm{\Gamma }_{10}\mathrm{\Gamma }_0`$ has an equal number of eigenvalues $`+1`$ and $`1`$. Thus we can chose a basis such that the only non zero brackets are
$`\{Q_a,Q_b\}=\delta _{ab},a,b=1,\mathrm{},16`$ (5.39)
and as in four dimensions half of the supercharges can be represented by zero. This is a general property of supersymmetric algebras . The non-zero supercharges are called the active supercharges. The representation of the eleven-dimensional supersymmetric algebra turns out to be the spinor representation of $`\text{Spin}(16)`$ (of dimension $`256`$), but as in four dimensions to identify the precise content of the multiplet we have to study the embedding $`\text{Spin}(16)\text{Spin}(9)`$, with $`\text{Spin}(9)`$ the little group. The active supercharges are in the $`\mathrm{𝟏𝟔}`$ (spinor) representation of $`\text{Spin}(9)`$ and in the $`\mathrm{𝟏𝟔}`$ (vector) representation of $`\text{Spin}(16)`$. To identify the representation of the supersymmetric algebra we proceed in several steps. We first observe that in the following embeddings we have the decomposition
$`\begin{array}{ccc}\text{Spin}(9)\hfill & \text{Spin}(8)\hfill & \text{Spin}(6)\times \text{Spin}(2)\hfill \\ \mathbf{16}\hfill & =\mathrm{𝟖}_+\mathrm{𝟖}_{}\hfill & =\left((\mathrm{𝟒}_+,\frac{1}{2})(\mathrm{𝟒}_{},\frac{1}{2})\right)\left((\mathrm{𝟒}_+,\frac{1}{2})(\mathrm{𝟒}_{},\frac{1}{2})\right)\hfill \end{array}`$ (5.42)
with $`\mathrm{𝟖}_\pm `$ a left/right handed spinor of $`\text{Spin}(8)`$, $`(\mathrm{𝟒}_\pm `$ a left/right handed spinor of $`\text{Spin}(6)`$ and $`\pm \frac{1}{2}`$ the eigenvalue of $`\text{Spin}(2)`$).
Then we study explicitly the spinor representation of $`\text{Spin}(8)`$. We denote $`Q_\pm =\mathrm{𝟖}_\pm `$, define a Clifford vacuum $`\mathrm{\Omega }_\pm `$ for each supercharge $`Q_\pm `$, and decompose along the line of remark 4.1 the $`Q_\pm `$ into operators of creation and annihilation (denoted $`(a_+,a_+^{})`$ and $`(a_{},a_{}^{})`$) and obtain the spinor representation (with $`Q_+`$ for instance). The supercharges $`Q_+`$ belong to the vector representation $`\mathrm{𝟖}_v`$ of some $`\text{Spin}(8)_Q`$ algebra generated by $`Q_+`$ -see Remark 4.1-, but they also belong to the spinor representation of the $`\text{Spin}(8)_{\text{s.t.}}`$ subgroup of the little group $`\text{Spin}(9)`$. Thus we firstly study the decomposition through the embedding $`\text{Spin}(8)_{\text{s.t.}}\text{Spin}(8)_Q`$ for which $`\mathrm{𝟖}_+=\mathrm{𝟖}_v`$
$`\begin{array}{ccc}\text{state}\hfill & \text{Spin}(6)\hfill & \text{Spin}(2)\hfill \\ \mathrm{\Omega }_+\hfill & \mathbf{1}\hfill & 1\hfill \\ a_+^{}\mathrm{\Omega }_+\hfill & \mathbf{4}_+\hfill & \frac{1}{2}\hfill \\ \left[a_+^{}\right]^2\mathrm{\Omega }_+\hfill & \mathbf{6}\hfill & \mathrm{\hspace{0.17em}0}\hfill \\ \left[a_+^{}\right]^3\mathrm{\Omega }_+\hfill & \mathbf{4}_{}\hfill & \frac{1}{2}\hfill \\ \left[a_+^{}\right]^4\mathrm{\Omega }_+\hfill & \mathbf{1}\hfill & \mathrm{\hspace{0.17em}1}\hfill \end{array}`$ (5.49)
where $`\left[a_+^{}\right]^n`$ means an $`n`$th antisymmetric product of operators of creation. In this decomposition, we have chosen the subgroup $`\text{Spin}(6)\times \text{Spin}(2)\text{Spin}(8)`$ such that $`a_+^{}=(\mathrm{𝟒}_+,\frac{1}{2})`$ this gives the second line in (5.49). Moreover, using $`(\mathrm{𝟒}_+,\frac{1}{2})(\mathrm{𝟒}_+,\frac{1}{2})=(\mathrm{𝟔},1)(\mathrm{𝟏𝟎}_+,1)`$ with $`\mathrm{𝟔}`$ the vector representation of $`\text{Spin}(6)`$ and $`\mathrm{𝟏𝟎}_+`$ the antisymmetric self-dual tensor of order three of $`\text{Spin}(6)`$, and oberving that $`\mathrm{𝟔}`$ corresponds to the antisymmetric tensor product of spinors and $`\mathrm{𝟏𝟎}_+`$ to the symmetric tensor product of spinors, this gives the third line in (5.49). Similar analysis give the $`\text{Spin}(6)`$ content of the other line of (5.49). Finally let us mention that the eigenvalues of $`\text{Spin}(2)`$ are normalized such that their sum is equal to zero. Now, it is easy to regroup the various terms and obtain representations of $`\text{Spin}(8)_{\text{s.t.}}`$:
$`\begin{array}{cc}\text{Spin}(6)\times \text{Spin}(2)\hfill & \text{Spin}(8)_{\text{s.t.}}\hfill \\ (\mathrm{𝟏},1)(\mathrm{𝟔},0)(\mathrm{𝟏},1)\hfill & =\mathrm{𝟖}_\text{v}\hfill \\ (\mathrm{𝟒}_+,\frac{1}{2})(\mathrm{𝟒}_{},\frac{1}{2})\hfill & =\mathrm{𝟖}_{}\hfill \end{array}`$ (5.53)
with $`\mathrm{𝟖}_\text{v}`$ the vector representation of $`\text{Spin}(8)_{\text{s.t.}}`$ and $`\mathrm{𝟖}_\pm `$ the two spinor representations. Thus in the embedding $`\text{Spin}(8)_{\text{s.t.}}\text{Spin}(8)_Q`$ we have the following decomposition $`\mathrm{𝟏𝟔}=\mathrm{𝟖}_{}\mathrm{𝟖}_\text{v}`$. In a similar way, acting with the supercharges $`Q_{}`$ we have the decomposition: $`\mathrm{𝟏𝟔}=\mathrm{𝟖}_+\mathrm{𝟖}_\text{v}`$. To obtain now the full $`\text{Spin}(16)`$ representation, we just have to tensorise the representations obtained with $`Q_+`$ and $`Q_{}`$. First we notice
$`\mathrm{𝟖}_\text{v}\mathrm{𝟖}_\text{v}`$ $`=`$ $`\mathrm{𝟏}[\mathrm{\Phi }]\mathrm{𝟑𝟓}[g_{ij}]\mathrm{𝟐𝟖}[B_{ij}]`$
$`\mathrm{𝟖}_+\mathrm{𝟖}_{}`$ $`=`$ $`\mathrm{𝟖}_\text{v}[A_i]\mathrm{𝟓𝟔}_\text{v}[C_{ijk}]`$ (5.54)
$`\mathrm{𝟖}_\text{v}\mathrm{𝟖}_+`$ $`=`$ $`\mathrm{𝟖}_{}[\lambda _R]\mathrm{𝟓𝟔}_{}[\mathrm{\Psi }_R{}_{}{}^{i}]`$
$`\mathrm{𝟖}_\text{v}\mathrm{𝟖}_{}`$ $`=`$ $`\mathrm{𝟖}_+[\lambda _L]\mathrm{𝟓𝟔}_+[\mathrm{\Psi }_L{}_{}{}^{i}]`$
with $`1i,j,k8`$ the $`\text{Spin}(8)_{\text{s.t.}}`$ indices and $`\mathrm{\Phi }`$ a scalar, $`g_{ij}`$ a symmetric traceless tensor, $`B_{ij}`$ a two-form, $`A_i`$ a vector, $`C_{ijk}`$ a three-form, $`\lambda _L`$ a left-handed spinor, $`\lambda _R`$ a right-handed spinor, $`\mathrm{\Psi }_L^i`$ a left-handed spinor-vector and $`\mathrm{\Psi }_R^i`$ a right-handed spinor-vector of $`\text{Spin}(8)_{\text{s.t.}}`$. The second decomposition comes from (4.76) The third and fourth decompositions come from the triality property of $`𝔰𝔬(8)`$ $``$look to the Dynkin diagram of $`𝔰𝔬(8)`$$``$. The field in bracket $`[]`$ just represents the type of field corresponding to the given representation. For instance $`\mathrm{𝟐𝟖}[B_{ij}]`$ means that the $`28`$dimensional representation of $`\text{Spin}(8)_{\text{s.t.}}`$ corresponds to a two-form $`B_{ij}`$. This gives the decomposition of the spinor representation of $`\text{Spin}(16)\text{Spin}(8)_{\text{s.t.}}`$. To obtain now the decomposition through the embedding $`\text{Spin}(16)\text{Spin}(9)`$, we just have to study the embedding $`\text{Spin}(9)\text{Spin}(8)_{\text{s.t.}}`$. If we define $`I,J,K=1,\mathrm{},9`$ the indices of $`\text{Spin}(9)`$ and $`i,j,k=1,\mathrm{},8`$ the indices of $`\text{Spin}(8)`$, we have:
$`\mathrm{𝟒𝟒}[g_{IJ}]`$ $`=`$ $`\mathrm{𝟑𝟓}[g_{ij}]\mathrm{𝟖}_\text{v}[g_{i9}]\mathrm{𝟏}[g_{99}]=\mathrm{𝟑𝟓}[g_{ij}]\mathrm{𝟖}_\text{v}[A_i]\mathrm{𝟏}[\mathrm{\Phi }]`$
$`\mathrm{𝟖𝟒}[C_{IJK}]`$ $`=`$ $`\mathrm{𝟓𝟔}_\text{v}[C_{ijk}]\mathrm{𝟐𝟖}[C_{ij9}]=\mathrm{𝟓𝟔}_\text{v}[C_{ijk}]\mathrm{𝟐𝟖}[B_{ij}]`$ (5.55)
$`\mathrm{𝟏𝟐𝟖}[\mathrm{\Psi }^I]`$ $`=`$ $`\mathrm{𝟓𝟔}_+[\mathrm{\Psi }_L^i]\mathrm{𝟖}_+[\mathrm{\Psi }_L^9]\mathrm{𝟓𝟔}_{}[\mathrm{\Psi }_R^i]\mathrm{𝟖}_{}[\mathrm{\Psi }_R^9]`$ (5.56)
$`=`$ $`\mathrm{𝟓𝟔}_+[\mathrm{\Psi }_L^i]\mathrm{𝟖}_+[\lambda _L]\mathrm{𝟓𝟔}_{}[\mathrm{\Psi }_R^i]\mathrm{𝟖}_{}[\lambda _R].`$
Thus finally the representation of the eleven-dimensional supersymmetric algebra contains one spinor-vector (a Rarita-Schwinger field), $`\mathrm{\Psi }^I`$, a symmetric traceless second order tensor (a tensor metric) and a three-form:
$`\mathrm{𝟐𝟓𝟔}=\mathrm{𝟏𝟐𝟖}_L\mathrm{𝟏𝟐𝟖}_R\{\begin{array}{cc}\mathrm{𝟏𝟐𝟖}_R=\mathrm{\Psi }^I\hfill & \text{fermion}\hfill \\ \mathrm{𝟏𝟐𝟖}_L=g_{IJ},C_{IJK}\hfill & \text{bosons}.\hfill \end{array}`$ (5.59)
These fields are representations of the little group $`\text{Spin}(9)`$. It becomes easy to obtain a representation of the Poincaré algebra not limiting the indices to their $`\text{Spin}(9)`$ values but allowing their $`\text{Spin}(1,10)`$ values. For instance $`g_{IJ},1I,J9g_{MN},0M,N10`$.
Compactifications
Having obtained the representation of the eleven-dimensional supersymmetric algebra we can now obtain the representations in smaller dimensions by compactification. Starting with the multiplet (5.59) the four dimensional representation of the extended $`N=8`$ supersymmetry can be built. In the little group $`SO(9)`$ the indices of the fields $`I,J=1,\mathrm{},9(i,j=1,2`$ and $`m,n=1,\mathrm{},7`$) corresponding to their $`\text{Spin}(2)`$ and $`\text{Spin}(7)`$ content. Thus we can decompose easily the fields in dimensional reduction from eleven to four dimensions. (In the simplest compactification, called the trivial compactification, we simply assume that the fields do not depend on the components of the compact dimension). This gives
$`\begin{array}{cccccc}g_{IJ}\hfill & \hfill & g_{ij}\hfill & g_{im}\hfill & g_{mn}\hfill & \\ & & 1\text{ graviton}\hfill & 7\text{ vectors}\hfill & 28\text{ scalars}\hfill & \\ C_{IJK}\hfill & \hfill & C_{ijk}\hfill & C_{ijm}\hfill & C_{imn}\hfill & C_{mnp}\hfill \\ & & \text{empty}\hfill & 7\text{ scalars}\hfill & 21\text{ vectors}\hfill & 35\text{ scalars}\hfill \\ \mathrm{\Psi }^I\hfill & \hfill & \mathrm{\Psi }^i\hfill & \mathrm{\Psi }^m\hfill & & \\ & & 8\text{ gravitinos}\hfill & 58\text{ spinors}\hfill & & \end{array}`$ (5.66)
In this decomposition we have to pay attention because with respect to $`\text{Spin}(2)`$ a three-form do not exists and a two-form is dual to a scalar. For the fermionic part of the multiplet we have to remember that through $`\text{Spin}(9)\text{Spin}(2)\times \text{Spin}(7)`$, we have the decomposition of a spinor $`\mathrm{𝟏𝟔}=(\frac{1}{2},\mathrm{𝟖})(\frac{1}{2},\mathrm{𝟖})`$. Finally, counting the number of states shows that we obtain the $`N=8`$ gravity multiplet.
###### Remark 5.6
The so-called Yang-Mills multiplet of the ten dimensional type I supersymmetry can be deduced from (5.53). The compactification of the type IIA supersymmetry from the eleven dimensional supersymmetry can be read off (5.3.2). In addition, using formulæ(5.3.2) with $`\mathrm{𝟖}_\pm \mathrm{𝟖}_\pm =\mathrm{𝟏}[\mathrm{\Phi }]\mathrm{𝟐𝟖}[B_{ij}]\mathrm{𝟑𝟓}[D_{ijkl}^\pm ]`$ with $`D_{ijkl}^\pm `$ an (anti-)self-dual four-form of $`\text{Spin}(8)`$ all ten dimensional supersymmetric multiplets can be obtained:
1. Yang-Mills multiplet type I: one vector $`A^i`$ one right-handed spinor $`\lambda _R`$. (The spinor is in the opposite chirality than the supercharge.)
2. Gravity type multiplet (the vacuum $`\mathrm{\Omega }`$ is a vector of $`\text{Spin}(8)`$): one scalar $`\mathrm{\Phi }`$, one tensor metric $`g_{ij}`$ one two-form $`B_{ij}`$; one left-handed spinor $`\lambda _L`$ one left-handed spinor-vector $`\mathrm{\Psi }_L^i`$. (The spinor and spinor-vector are in the same chirality than the supercharge).
3. type IIA one scalar $`\mathrm{\Phi }`$, one tensor metric $`g_{ij}`$ one two-form $`B_{ij}`$, one vector $`A_i`$, one three-form $`C_{ijk}`$, one left-handed and one right-handed spinor $`\lambda _R,\lambda _L`$ and one left-handed and one right-handed spinor-vector $`\mathrm{\Psi }_L^i,\mathrm{\Psi }_R^i`$. The fermions are of both chirality.
4. type IIB: two scalars $`\mathrm{\Phi },\mathrm{\Phi }^{}`$, one tensor metric $`g_{ij}`$, two two-forms $`B_{ij},B_{ij}^{}`$ one anti-self-dual four form $`D_{ijkl}^{}`$ two left handed spinors $`\lambda _L,\lambda _L^{}`$ two left-handed spinor-vector $`\mathrm{\Psi }_L^i,\mathrm{\Psi }_L^i`$.
Finally let us mention that type IIA or type IIB supersymmetry give by compactification in four dimensions $`N=8`$ extended-supersymmetry.
## 6 Conclusion
In this lecture we have shown, that the basic tools to construct supersymmetric extensions of the Poincaré algebra are Clifford algebras. Special attention have been given to the four, ten and eleven dimensional spaces-times. Studying representations of supersymmetric algebras shows that a supermultiplet contains an equal number of bosonic and fermionic degrees of freedom and that supersymmetry is a symmetry which mixes non-trivially bosons and fermions. The next step is to apply supersymmetric theories in Quantum Field Theory or particle physics. For that purpose we need first to calculate transformation of the fields under supersymmetry and then to build invariant Lagrangians (the concept of superspace is central for this construction). For instance all the technics of supersymmetry have been applied in four dimensions to construct a supersymmetric version of the standard model . There are some strong arguments in favor of such a theory, even if there is no experimental evidence of supersymmetry. Supersymmetry or more precisely its local version contains gravity and as such is named supergravity. Ten-dimensional supergravities appear as some low energy limits of string theories and present some interesting duality properties . Finally, with the brane revolution a lot of hope have been put to the so-called M-theory whose low energy limits contains the eleven-dimensional supergravity and the various strings theory .
##### Acknowledgements
The Organizing Committee, and in particular P. Angles, are kindly acknowledged for the friendly and studious atmosphere during the conference. |
warning/0506/math0506111.html | ar5iv | text | # Orbifold Quantum Riemann-Roch, Lefschetz and Serre
## 1. Introduction
Our main goal is to extend the Quantum Riemann-Roch theorem of Coates-Givental in Gromov-Witten theory to the case of algebraic orbifold target spaces (i.e. smooth Deligne-Mumford stacks). As applications, we prove Quantum Serre duality for Deligne-Mumford stacks and a general form of Quantum Lefschetz Hyperplane Theorem for Deligne-Mumford stacks. Results leading towards mirror symmetry of complete intersection orbifolds are also deduced as consequences.
### 1.1. Background: Gromov-Witten Theory of Stacks
We work over the field of complex numbers $``$. A Deligne-Mumford stack $`𝒳`$ is a category fibered in groupoids which satisfies several rather complicated conditions. For the precise definition and detailed discussions about properties of Deligne-Mumford stacks, we refer to and . It is known that a (separated) Deligne-Mumford stack $`𝒳`$ has a coarse moduli space $`X`$ which is in general an algebraic space. For any closed point $`xX`$ there is an étale neighborhood $`U_xX`$ of $`x`$ such that the pullback $`U_x\times _X𝒳`$ is a stack of the form $`[V_x/\mathrm{\Gamma }_x]`$ with $`V_x`$ affine and $`\mathrm{\Gamma }_x`$ a finite group. Thus one may view a Deligne-Mumford stack as a geometric object locally a quotient of an affine scheme by a finite group, just like one would view a scheme as a geometric object locally an affine scheme. This viewpoint is in analogy with the notion of orbifolds in differential geometry: A complex orbifold is a topological space $`X`$ together with a choice of an open neighborhood $`U_xx`$ for each $`xX`$, an open subset $`V_x^D`$, and a finite group $`\mathrm{\Gamma }_x`$ acting on $`V_x`$ such that $`U_x`$ is homeomorphic to a quotient $`V_x/\mathrm{\Gamma }_x`$ of $`V_x`$ by a finite group $`\mathrm{\Gamma }_x`$ and the collection $`\{U_x,V_x,\mathrm{\Gamma }_x\}_{xX}`$ satisfies some compatibility conditions concerning $`\mathrm{\Gamma }_x`$-actions on overlaps.
In this paper we work with Deligne-Mumford stacks, but in view of the analogy mentioned above, the term “orbifold” will also be used. By abuse of language, we will treat the terms “orbifold” and “smooth Deligne-Mumford stack” as synonymous<sup>1</sup><sup>1</sup>1We do not assume that a Deligne-Mumford stack has trivial generic stabilizers, unless otherwise mentioned..
Let $`𝒳`$ be a smooth Deligne-Mumford stack with projective coarse moduli space $`X`$. The inertia stack $`I𝒳:=𝒳\times _{\mathrm{\Delta },𝒳\times 𝒳,\mathrm{\Delta }}𝒳`$ associated to $`𝒳`$ plays an important role in the theory of stacks. Locally at $`x𝒳`$, the inertia stack $`I𝒳`$ consists of connected components labeled by conjugacy classes of elements $`g\mathrm{\Gamma }_x`$. Each connected component is described locally as a quotient $`[V_x^g/C_{\mathrm{\Gamma }_x}(g)]`$, where $`V_x^gV_x`$ denotes the locus fixed by $`g`$ and $`C_{\mathrm{\Gamma }_x}(g)\mathrm{\Gamma }_x`$ denotes the centralizer of $`g`$. Objects in the category underlying $`I𝒳`$ are pairs $`(x,g)`$ with $`x`$ an object in $`𝒳`$ and $`gAut_𝒳(x)`$. There is a canonical projection $`q`$ from $`I𝒳`$ to $`𝒳`$. Also, $`I𝒳`$ contains $`𝒳`$ as the component corresponding to choosing $`g`$ to be the identity element in $`\mathrm{\Gamma }_x`$. See Section 2.1 for more details.
The construction of Gromov-Witten invariants as intersection numbers on the moduli spaces of stables maps was generalized to symplectic orbifolds by Chen-Ruan and to Deligne-Mumford stacks by Abramovich-Graber-Vistoli , . A summary of the basics of Gromov-Witten theory for stacks will be given in Section 2. The ideas central to their constructions are: (1) the domain curves $`𝒞`$ of a stable map $`𝒞𝒳`$ to a stack can be orbicurves, i.e. they can have nontrivial stack structures at marked points and nodes; (2) the stable maps $`𝒞𝒳`$ are required to respect the stack structures of $`𝒞`$ and $`𝒳`$, i.e. they should be representable morphisms.
In this paper, we consider a variant of Gromov-Witten theory for stacks. Suppose that $`𝒳`$ satisfies Assumption 2.5.9 below. Given a complex vector bundle $`F`$ on $`𝒳`$ and an invertible multiplicative characteristic class $`𝐜()`$ of complex vector bundles, we define twisted orbifold Gromov-Witten invariants using these data. These twisted invariants can be encoded in a generating function, called $`(𝐜,F)`$-twisted total descendant potential of $`𝒳`$, which is defined as follows:
$$𝒟_{(𝐜,F)}(𝐭):=\mathrm{exp}\left(\underset{g=0}{\overset{\mathrm{}}{}}\mathrm{}^{g1}\underset{n,d}{}\frac{Q^d}{n!}_{[\overline{}_{g,n}(𝒳,d)]^w}𝐜(F_{g,n,d})_{i=1}^n\underset{k=0}{\overset{\mathrm{}}{}}\mathrm{ev}_i^{}(t_k)\overline{\psi }_i^k\right).$$
Let us explain the notations in this definition. Integration in this formula is performed over the weighted virtual fundamental class $`[\overline{}_{g,n}(𝒳,d)]^w`$ in the moduli space $`\overline{}_{g,n}(𝒳,d)`$ of degree-$`d`$ stable maps to $`𝒳`$ from genus-$`g`$ orbicurves with sections to all $`n`$ marked gerbes. The cohomology classes $`t_kH^{}(I𝒳,)`$ for $`k=0,1,2,\mathrm{}`$, are pulled back to the moduli space by the evaluation maps<sup>2</sup><sup>2</sup>2Due to presence of stack structures on the domain curves, the evaluation of stable maps at marked points takes values in $`I𝒳`$ (rather then $`𝒳`$). Note that $`𝒳`$ is a component of $`I𝒳`$. $`\mathrm{ev}_i:\overline{}_{g,n}(𝒳,d)I𝒳`$, $`i=1,\mathrm{},n`$. The classes $`\overline{\psi }_i`$ are the first Chern classes of the universal cotangent line<sup>3</sup><sup>3</sup>3These are the cotangent line of the underlying coarse curve, not the orbicurve. See Section 2.5.1 for details. bundles over the moduli spaces $`\overline{}_{g,n}(𝒳,d)`$. The “twisting factor” $`𝐜(F_{g,n,d})`$ is the characteristic class $`𝐜`$ applied to the virtual bundle $`F_{g,n,d}K^0(\overline{}_{g,n}(𝒳,d))`$, which is constructed as follows: Consider the universal family of orbifold stable maps,
$$\begin{array}{ccc}𝒞_{g,n}(𝒳,d)& \stackrel{\mathrm{ev}}{}& 𝒳\\ f& & \\ \overline{}_{g,n}(𝒳,d).\end{array}$$
By definition, $`f`$ is a family of nodal orbicurves which are the source curves of the orbifold stable maps, and the restrictions of $`\mathrm{ev}`$ to the fibers give rise the stable maps which $`\overline{}_{g,n}(𝒳,d)`$ parametrizes. We put<sup>4</sup><sup>4</sup>4It follows from the results of that the map $`f`$ is a local complete intersection morphism. Therefore the K-theoretic push-forward $`Rf_{}`$ of a bundle has a locally free resolution and thus defines an element in the Grothendieck group $`K^0`$. See Appendix B for more discussion on this.
$$F_{g,n,d}:=Rf_{}(\mathrm{ev}^{}F)K^0(\overline{}_{g,n}(𝒳,d)).$$
$`Q^d`$ is an element in the Novikov ring $`\mathrm{\Lambda }_{nov}`$ (see Section 2.5.2), and $`\mathrm{}`$ is a formal variable. Finally, $`𝒟_{(𝐜,F)}(𝐭)`$ depends on $`(t_0,t_1,t_2,\mathrm{})`$ and we package them as $`𝐭(z)=_{k0}t_kz^k`$. (Although we denote it by $`𝒟_{(𝐜,F)}(𝐭)`$, the descendant potential does not depend on $`z`$.)
The “untwisted” total descendant potential $`𝒟_𝒳`$ of $`𝒳`$, which encodes usual orbifold Gromov-Witten invariants, is defined by the defining equation of $`𝒟_{(𝐜,F)}`$ with the twisting factors $`𝐜(F_{g,n,d})`$ replaced by $`1`$. Details of the definition of twisted orbifold Gromov-Witten invariants will be given in Section 2.5.8.
### 1.2. Main Result: Orbifold Quantum Riemann-Roch
The main result of this paper, orbifold Quantum Riemann-Roch theorem, expresses the twisted orbifold Gromov-Witten invariants in terms of the usual invariants. To state the result we need the following quantization formalism introduced into Gromov-Witten theory in . Here we give a brief summary, see Section 3 for a detailed treatment. Let $`H:=H^{}(I𝒳,)`$ be the cohomology (super-)space of the inertia stack. The space $`H`$ is equipped with the symmetric inner product (called the orbifold Poincaré pairing)
$$(a,b)_{orb}:=_{I𝒳}aI^{}b,a,bH,$$
where $`I`$ is an involution on $`I𝒳`$ induced by the inversion $`gg^1,gAut_𝒳(x),x𝒳`$. Fix an additive basis $`\{\varphi _\alpha \}`$ of $`H`$ and let $`\{\varphi ^\alpha \}`$ be the dual basis with respect to $`(,)_{orb}`$. Introduce the space $`:=H\mathrm{\Lambda }_s\{z,z^1\}`$ of convergent Laurent series in $`z`$ (see Section 3.1). Following and , we equip $``$ with the $`\mathrm{\Lambda }_s`$-valued even symplectic form
$$\mathrm{\Omega }(f,g):=\text{Res}_{z=0}(f(z),g(z))_{orb}dz,f,g.$$
The Lagrangian polarization $`=_+_{}`$, with $`_+=H\mathrm{\Lambda }_s\{z\}`$ and $`_{}=z^1(H\mathrm{\Lambda }_s\{z^1\})`$, identifies $`(,\mathrm{\Omega })`$ with the cotangent bundle $`T^{}_+`$, see Section 3.1. Let $`p_a^\mu ,q_b^\nu `$ be Darboux coordinates of $`(,\mathrm{\Omega })`$ with respect to this polarization, as introduced in Section 3.1. Put $`p_k:=_\mu p_k^\mu \varphi ^\mu ,q_k=_\nu q_k^\nu \varphi _\nu `$ and $`(𝐩:=_{k0}p_k(z)^{k1},𝐪:=_{k0}q_kz^k)`$.
Let $`ch_k()`$ denote the degree $`2k`$ component of the Chern character. We may view $`𝐜()=\mathrm{exp}(_ks_kch_k())`$ as a family of characteristic classes depending on variables $`s=(s_0,s_1,\mathrm{})`$. As $`s`$ varies, the twisted descendent potentials $`𝒟_{(𝐜,F)}`$ define a family $`𝒟_s`$ of elements in the Fock space<sup>5</sup><sup>5</sup>5See Section 3.1 for definition. of formal functions on $`_+`$ using the following convention: For $`𝐭(z)=t_0+t_1z+t_2z^2+\mathrm{}H\mathrm{\Lambda }_s\{z\}`$, we identify $`𝐭`$ with the Darboux coordinates $`𝐪H\mathrm{\Lambda }_s\{z\}`$ via
$$𝐪(z)=\sqrt{𝐜(F^{(0)})}(𝐭(z)\mathrm{𝟏}z),$$
where $`F^{(0)}`$ is the vector bundle on $`I𝒳`$ whose fiber at $`(x,g)I𝒳`$ is the subspace of $`F|_x`$ on which $`g`$ acts with eigenvalue $`1`$, and $`\mathrm{𝟏}H^{}(I𝒳,)`$ is the unit cohomology class of the principal component $`𝒳I𝒳`$ (see Section 2.1). Then put
$$𝒟_s(𝐪):=𝒟_{(𝐜,F)}(𝐭).$$
In other words, $`𝒟_{(𝐜,F)}`$ is now viewed as a function in $`q_0=\sqrt{𝐜(F^{(0)})}t_0,q_1=\sqrt{𝐜(F^{(0)})}(t_0\mathrm{𝟏})`$ and $`q_k=\sqrt{𝐜(F^{(0)})}t_k,k2`$. Note that $`𝒟_s|_{s_0=s_1=\mathrm{}=0}=𝒟_𝒳`$.
We also need to define certain vector bundles on the inertia stack. The inertia stack is a disjoint union
$$I𝒳=\underset{i}{}𝒳_i,$$
where $``$ is a index set. For any $`(x,g)𝒳_i`$ let $`r_i`$ denote the order of the element $`gAut_𝒳(x)`$. To a vector bundle $`F`$ on $`𝒳`$, define $`F_i^{(l)}`$ over $`𝒳_i`$ to be the vector bundle whose fiber $`F_i^{(l)}|_{(x,g)}`$ at $`(x,g)𝒳_i`$ is the subspace of $`F|_x`$ on which $`g`$ acts with eigenvalue $`\mathrm{exp}(2\pi \sqrt{1}l/r_i)`$. See Section 2.2 for more details. Also observe that $`H^{}(I𝒳,)=_iH^{}(𝒳_i,)`$.
The following is the main result of this paper.
###### Theorem 1 (Orbifold Quantum Riemann-Roch, see Theorem 4.2.1).
$$𝒟_s\widehat{\mathrm{\Delta }}𝒟_𝒳.$$
Here $`\mathrm{\Delta }:`$ is the operator given by ordinary multiplication by
$$\mathrm{\Delta }=\sqrt{𝐜(F^{(0)})}\underset{i}{}\mathrm{exp}\left(\underset{0lr_i1}{}\underset{k0}{}s_k\underset{m0}{}\frac{B_m(l/r_i)}{m!}\mathrm{ch}_{k+1m}(F_i^{(l)})z^{m1}\right),$$
and $`\widehat{\mathrm{\Delta }}`$ is the differential operator obtained by quantizing $`\mathrm{\Delta }`$.
We now explain the ingredients in the Theorem.
1. The symbol $``$ stands for “equal up to a scalar factor depending on $`s`$” which will be explicitly described in Section 4; see Theorem 4.2.1.
2. Here $`B_m(x)`$ are the Bernoulli polynomials defined by
$$\frac{te^{tx}}{e^t1}=\underset{m0}{}\frac{B_m(x)t^m}{m!}.$$
For example, $`B_0(x)=1,B_1(x)=x1/2,B_2(x)=x^2x+1/6`$.
3. The operators on $``$ defined as multiplication by $`\mathrm{ch}_{k+1m}(F_i^{(l)})z^{m1}`$ over the component $`𝒳_i`$ of $`I𝒳`$ turns out to be anti-symmetric with respect to the form $`\mathrm{\Omega }`$ and thus define infinitesimal linear symplectic transformations on $``$, see Corollary 4.1.5. The quantized operator $`\widehat{\mathrm{\Delta }}`$ on the Fock space is defined as follows: The operator
$$\mathrm{log}\mathrm{\Delta }:=\frac{1}{2}\underset{k0}{}s_k\mathrm{ch}_k(F^{(0)})+\underset{i}{}\underset{0lr_i1}{}\underset{k0}{}s_k\underset{m0}{}\frac{B_m(l/r_i)}{m!}\mathrm{ch}_{k+1m}(F_i^{(l)})z^{m1}$$
is infinitesimally symplectic. We define $`\widehat{\mathrm{\Delta }}:=\mathrm{exp}(\widehat{\mathrm{log}\mathrm{\Delta }})`$, where $`\widehat{\mathrm{log}\mathrm{\Delta }}`$ is the differential operator defined by quantizing the quadratic Hamiltonians of $`\mathrm{log}\mathrm{\Delta }`$ following the standard rule in Darboux coordinates:
$$(q_\alpha q_\beta )\widehat{}:=\mathrm{}^1q_\alpha q_\beta ,(p_\alpha p_\beta )\widehat{}:=\mathrm{}_{q_\alpha }_{q_\beta },(q_\alpha p_\beta )\widehat{}:=q_\alpha _{q_\beta }.$$
See Section 3.3 for more details on the quantization procedure.
###### Remark 1.
1. When the target space $`𝒳`$ is a manifold, $`\mathrm{\Delta }`$ is simplified to
$$\mathrm{exp}\left(\underset{k0}{}s_k\underset{m0}{}\frac{B_{2m}(0)}{(2m)!}\mathrm{ch}_{k+12m}(F)z^{2m1}\right).$$
Thus our main Theorem recovers the Quantum Riemann-Roch theorem of Coates-Givental . Their proof is based on the Grothendieck-Riemann-Roch (GRR) theorem applied to a family of nodal curves and thus goes back to Mumford and Faber-Pandharipande . Our proof of Theorem 1 relies on an appropriate generalization, in the spirit of Kawasaki , of the GRR formula valid for morphisms between Deligne-Mumford stacks. This version of the GRR formula, explained in Appendix A, is known to hold in algebraic context (it is a result of B. Toen). It is tempting to extend our results to almost Kähler orbifolds, but we are unable to do so at this moment. The case of almost Kähler manifolds is treated in Appendix B of .
2. The Bernoulli numbers $`B_{2m}(0)`$ arise naturally in the formula of Coates-Givental due to the use of the GRR formula. Peculiarly, the values $`B_m(l/r)`$ of the Bernoulli polynomials featuring in our main result do not seem to arise in the generalization of the GRR formula to the case of orbifolds. It would be interesting to have a conceptual understanding of the presence of Bernoulli polynomials in our result.
Theorem 1 has some immediate consequences in genus zero. The genus zero $`(𝐜,F)`$-twisted descendant potential is defined as
$$_{(𝐜,F)}^0:=\underset{n,d}{}\frac{Q^d}{n!}_{[\overline{}_{0,n}(𝒳,d)]^w}𝐜(F_{0,n,d})_{i=1}^n\underset{k=0}{\overset{\mathrm{}}{}}\mathrm{ev}_i^{}(t_k)\overline{\psi }_i^k.$$
It is viewed as a element in the Fock space in the way described above. The genus zero orbifold Gromov-Witten potential $`_𝒳^0`$ is defined by the above equation with the twisting factor $`𝐜(F_{0,n,d})`$ replaced by $`1`$. The graphs of the differentials of $`_{(𝐜,F)}^0`$ and $`_𝒳^0`$ are two (formal germs of) Lagrangian submanifolds $`_s=_{(𝐜,F)}`$ and $`_𝒳`$ of the symplectic vector space $``$. Theorem 1 yields a relationship between these two Lagrangian submanifold germs.
###### Corollary 1 (=Corollary 4.2.3).
$$_s=\mathrm{\Delta }_𝒳.$$
The genus $`0`$ orbifold Gromov-Witten potential $`_𝒳^0`$ is known to satisfy three sets of partial differential equations: the string equation (SE), the dilaton equation (DE), and topological recursion relations (TRR), see Section 2.5.7. According to Givental (see , Theorem 1), this is equivalent to the following property of the Lagrangian submanifold germ $`_𝒳`$:
Property ($``$): $`_𝒳`$ is the germ of a Lagrangian cone with the vertex at the origin and such that its tangent spaces $`L`$ are tangent to $`_𝒳`$ exactly along $`zL`$. See (3.1.1.1) for its precise meaning.
The property ($``$) is formulated in terms of the symplectic structure $`\mathrm{\Omega }`$ and the operator of multiplication by $`z`$. It does not depend on the choice of polarization. Therefore it is invariant under the action of the twisted loop group, which consists of $`End(H^{}(I𝒳))`$-valued formal Laurent series $`M`$ in $`z^1`$ satisfying<sup>6</sup><sup>6</sup>6Here $``$ denotes the adjoint with respect to $`(,)_{orb}`$. $`M^{}(z)M(z)=1`$. One checks that $`\mathrm{\Delta }`$ defines an element in the twisted loop group. This yields the following corollary:
###### Corollary 2.
The Lagrangian submanifold $`_s`$ satisfies property ($``$). In other words, twisted orbifold Gromov-Witten invariants in genus zero satisfy the axioms (TRR), (SE), and (DE) of genus zero theory.
### 1.3. Applications of Quantum Riemann-Roch
#### 1.3.1. Quantum Serre duality
Consider the orbifold Gromov-Witten theory twisted by the dual vector bundle $`F^{}`$ and the dual class
$$𝐜^{}():=\mathrm{exp}\left(\underset{k0}{}(1)^{k+1}s_kch_k()\right).$$
Theorem 1 implies the following “Quantum Serre duality”.
###### Corollary 3 (=Theorem 6.1.1).
Let $`𝐭^{}(z)=𝐜(F)𝐭(z)+(\mathrm{𝟏}𝐜(F))z`$. Then we have
$$𝒟_{(𝐜^{},F^{})}(𝐭^{})𝒟_{(𝐜,F)}(𝐭).$$
See Theorem 6.1.1 for the precise $`s`$-dependent scalar factor.
We may equip the bundle $`F`$ with an $`^{}`$-action given by scaling the fibers. We are also interested in the special case of twisting by the equivariant Euler class $`e()`$ with respect to this $`^{}`$-action. Let the dual bundle $`F^{}`$ be equipped with the dual $`^{}`$-action and let $`e^1()`$ be the inverse $`^{}`$-equivariant Euler class.
Let $`M:H^{}(I𝒳)H^{}(I𝒳)`$ be the operator defined as follows: On the cohomology $`H^{}(𝒳_i)`$ of a component $`𝒳_iI𝒳`$, $`M`$ is defined to be multiplication by the number $`(1)^{age(F_i)+\frac{1}{2}rankF_i^{mov}}`$. See Section 6.2 for more details, including definitions of $`age(F_i)`$, $`(q^{}F)^{inv}`$, and $`F_i^{mov}`$.
Put $`𝐭^{}(z)=z+(1)^{\frac{1}{2}rank(q^{}F)^{inv}}Me(F)(𝐭(z)\mathrm{𝟏}z)`$ and define a change
$$:Q^dQ^d(1)^{c_1(F),d},Q^d\mathrm{\Lambda }_{nov},$$
in the Novikov ring. The following Proposition is deduced from Corollary 3.
###### Proposition 1 (=Theorem 6.2.1).
$$𝒟_{(e^1,F^{})}(𝐭^{},Q)𝒟_{(e,F)}(𝐭,Q).$$
See Theorem 6.2.1 for the precise constant factor.
#### 1.3.2. Quantum Lefschetz
Again we consider the $`^{}`$-action on the bundle $`F`$ given by scaling the fibers. Let $`\lambda `$ denote the equivariant parameter. We now consider the genus zero theory of the special case of twisting by $`^{}`$-equivariant Euler class $`e`$ of this action. We assume that $`F`$ is pulled back from the coarse moduli space $`X`$. In this situation, the operator $`\mathrm{\Delta }`$ is closely related to asymptotics of the Gamma function:
$$\mathrm{\Delta }\frac{1}{\sqrt{e(F)}}\underset{i=1}{\overset{N}{}}\frac{1}{\sqrt{2\pi z}}_0^{\mathrm{}}e^{\frac{x+(\lambda +q^{}\rho _i)\mathrm{ln}x}{z}}𝑑x,$$
where $`\rho _1,\mathrm{},\rho _N`$ are Chern roots of $`F`$.
The intersection of $`_𝒳`$ with the affine subspace $`z+z_{}`$ defines a function $`J_𝒳(t,z)`$ called the $`J`$-function: For $`tH^{}(I𝒳)`$, define
$$J_𝒳(t,z):=z+t+d_𝐪_𝒳^0|_{𝐪=tz},$$
see Definition 3.1.2 for more explanation.
###### Theorem 2 (=Theorem 5.1.6).
Let $`F`$ be a vector bundle which is a direct sum of line bundles pulled back from the coarse moduli space $`X`$. Let $`\rho _1,\mathrm{},\rho _N`$ be the Chern roots of $`F`$. Let a formal function $`I(t,z)`$ of $`tH`$ be given as in Definition 5.1.4. Then the family
$$tI(t,z),tH$$
lies on the cone $`_{(e,F)}`$. In view of the property ($``$), the cone $`_{(e,F)}`$ is determined by this family.
This is an abstract form of the Quantum Lefschetz Hyperplane Theorem for Deligne-Mumford stacks, which generalizes previous results for varieties (see , , , , , ).
###### Remark 1.3.3.
The formal function $`I(t,z)`$ may also be written as follows:
$$I(t,z)=\underset{i=1}{\overset{N}{}}\frac{_0^{\mathrm{}}e^{x/z}J_𝒳(t+(\lambda +\rho _i)\mathrm{ln}x,z)𝑑x}{_0^{\mathrm{}}e^{\frac{x(\lambda +q^{}\rho _i)\mathrm{ln}x}{z}}𝑑x},$$
where the integrals are interpreted as their stationary phase asymptotics as $`z0`$. To see this, first rewrite $`J_𝒳`$ in the above equation using the string and divisor equations, then use integration by parts.
### 1.4. Towards a mirror theorem for orbifolds
Many examples of Calabi-Yau varieties in the mathematics and physics literatures are constructed as complete intersections in toric varieties, and many of them have quotient singularities. In dimension at most three, one may avoid dealing with singular Calabi-Yaus by taking crepant resolutions. In higher dimension, this is not possible in general since crepant resolutions may not exist. Therefore it is desirable to work directly with varieties with quotient singularities. The structure of quotient singularities on varieties is naturally described via Deligne-Mumford stacks.
A motivation to introduce twisted Gromov-Witten invariants is to compute Gromov-Witten invariants of complete intersections and verify predictions from mirror symmetry of Calabi-Yau manifolds (for example quintic threefolds in $`^4`$). This approach first appeared in the work of Kontsevich . Ever since the formulation of Quantum Lefschetz hyperplane principle (see e.g. , , ), the verification of mirror symmetry predictions for complete intersections has been divided into two independent parts:
1. Compute Gromov-Witten invariants for the ambient spaces;
2. Understand relationships between Gromov-Witten invariants of the complete intersections and those of the ambient spaces.
One motivation of the present paper is to prove mirror symmetry predictions for orbifolds using this approach. The Quantum Lefschetz hyperplane theorem proved in this paper establishes part (2) for orbifold target spaces under additional assumptions. A more useful version of quantum Lefschetz theorem for orbifolds is proven in . So far, works on part (1) have been most successful in the case of toric varieties. The toric mirror construction (see for instance ) applied to a toric orbifold $`𝒳`$ yields conjectural mirror pairs of Calabi-Yau orbifolds as complete intersections in toric orbifolds. Under additional convexity assumptions, some twisted orbifold Gromov-Witten invariants are related to orbifold Gromov-Witten invariants of the complete intersections. Thus our Quantum Lefschetz Hyperplane Theorem gives relations between genus-0 orbifold Gromov-Witten invariants of those Calabi-Yau orbifolds and the invariants of the ambient toric orbifolds, see Corollary 5.2.6. Once the orbifold Gromov-Witten invariants of toric orbifolds are computed (i.e. part (1) is settled), our result yields information about genus-0 orbifold Gromov-Witten invariants of the Calabi-Yau complete intersection orbifolds. This will eventually lead to verifications of mirror symmetry prediction for toric complete intersection orbifolds. Using the results of , the case of complete intersections in weighted projective spaces is treated in . We hope to return to other cases in the future.
### 1.5. Plan of the paper
The rest of the paper is organized as follows. Sections 2 and 3 contain most of the preparatory materials. In Section 2 we present some definitions and properties used throughout this paper. Section 2.1 and 2.2 contain discussions on important notions of stacks needed in this paper. Properties of orbifold cohomology are reviewed in Section 2.3. Section 2.4 is devoted to the moduli spaces of orbifold stable maps, on which orbifold Gromov-Witten theory is based. In Section 2.5 we review the orbifold Gromov-Witten theory constructed in and . We introduce the twisted orbifold Gromov-Witten invariants in Section 2.5.8. In Section 3 we explain how Givental’s symplectic vector space formalism , can be applied to twisted and untwisted orbifold Gromov-Witten theory. In Section 4 we state the orbifold Quantum Riemann-Roch theorem (Theorem 4.2.1). This is used to derive Quantum Lefschetz Hyperplane Principle in Section 5.1 and 5.2. Orbifold Quantum Serre duality is proved in Section 6. Section 7 contains a proof of Theorem 4.2.1. We discuss a Grothendieck-Riemann-Roch formula for Deligne-Mumford stacks in Appendix A. Appendix B concerns properties of the virtual bundle $`F_{g,n,d}`$. Some calculation concerning the quantized operators are given in Appendices C and D. In Appendix E we present a proof of the topological recursion relation for genus $`0`$ orbifold Gromov-Witten theory.
## Acknowledgments
The author is deeply grateful to A. Givental for his guidance, constant help and encouragement. Many thanks to D. Abramovich, T. Coates, T. Graber, and H. Iritani for numerous helpful discussions and suggestions on the subject and their interests in this work. Thanks to A. Kresch, Y.-P. Lee, M. Olsson, and B. Toen for many helpful discussions. The author is grateful for the referees’ numerous helpful comments and suggestions, which greatly improved the paper. The first version of this paper forms the main part of the author’s Ph.D. thesis. During the subsequent revision of this paper, the author was supported in part by postdoctoral fellowships from the Mathematical Science Research Institute (Berkeley, California) and Pacific Institute of Mathematical Sciences (Vancouver, Canada), and a visiting fellowship from Institut Mittag-Leffler (Djursholm, Sweden).
## 2. Orbifolds and their Gromov-Witten theory
In this section, we present some definitions, notations and properties which we use throughout.
### 2.1. Orbifolds
Throughout this paper, let $`𝒳`$ be a proper smooth Deligne-Mumford stack over the complex numbers $``$ with projective coarse moduli space $`X`$. In this section, we discuss some general properties of $`𝒳`$ and fix notations throughout.
A friendly introduction to basic notions of stacks can be found in . For comprehensive introductions to rigorous foundation of stacks the reader may consult and the Appendix of . A very detailed treatment of the theory of algebraic stacks can be found in (see also the forthcoming book ). The geometry of a stack of the form $`[M/G]`$ with $`M`$ a scheme and $`G`$ an algebraic group is essentially equivalent to the equivariant geometry of $`M`$ with respect to the $`G`$-action. Since almost all stacks we treat in this paper are of this form, keeping this interpretation in mind may help the readers unfamiliar with stacks understand this paper.
Recall that a morphism $`f:𝒳𝒴`$ of stacks is called representable if for every morphism $`g:S𝒴`$ from a scheme $`S`$, the fiber product $`S\times _{g,𝒴,f}𝒳`$ is a scheme. In particular, any morphism from a scheme to a stack is representable.
To a Deligne-Mumford stack $`𝒳`$ we can associate a coarse moduli space $`X`$ which is in general an algebraic space . For a morphism $`𝒳𝒴`$ of stacks, there is an induced morphism $`XY`$ between their coarse moduli spaces.
We now introduce the inertia stack associated to a stack $`𝒳`$, which plays a central role in Gromov-Witten theory for stacks.
###### Definition 2.1.1.
Let $`𝒳`$ be a Deligne-Mumford stack. The inertia stack $`I𝒳`$ associated to $`𝒳`$ is defined to be the fiber product
$$I𝒳:=𝒳\times _{\mathrm{\Delta },𝒳\times 𝒳,\mathrm{\Delta }}𝒳,$$
where $`\mathrm{\Delta }:𝒳𝒳\times 𝒳`$ is the diagonal morphism. The objects in the category underlying $`I𝒳`$ can be described as follows:
$$\begin{array}{cc}\hfill Ob(I𝒳)& =\{(x,g)|xOb(𝒳),gAut_𝒳(x)\}\hfill \\ & =\{(x,H,g)|xOb(𝒳),HAut_𝒳(x),g\text{ a generator of }H\}.\hfill \end{array}$$
###### Remark 2.1.2.
1. For a stack $`𝒳`$ over $``$, $`I𝒳`$ is isomorphic to the stack of representable morphisms from a constant cyclotomic gerbe to $`𝒳`$,
(2.1.2.1)
$$I𝒳\underset{r}{}HomRep(B\mu _r,𝒳).$$
At the level of objects, this means
$$Ob(I𝒳)=\{(x,H,\chi )|xOb(𝒳),HAut(x),\chi :H\mu _r\text{ an isomorphism for some }r\}.$$
Since we work over $``$ we will from now identify $`\mu _r`$ as the subgroup of $`^{}`$ of $`r`$-th roots of $`1`$, and fix a generator $`𝔲_r:=\mathrm{exp}(2\pi \sqrt{1}/r)`$ of $`\mu _r`$. In doing so, the identification (2.1.2.1) can be described as follows. An object $`(x,g)`$ of $`I𝒳`$ over a scheme $`S`$ is identified with a representable morphism $`S\times B\mu _r𝒳`$ such that the image is $`x`$ and the induced group homomorphism $`\mu _rAut_𝒳(x)`$ takes $`𝔲_r`$ to $`g`$.
This description of $`I𝒳`$ will also be used. For more details, see , Section 4.4 and , Section 3.2.
2. There is a natural projection $`q:I𝒳𝒳`$. On objects we have $`q((x,g))=x`$.
An important observation is that the inertia stack $`I𝒳`$ is in general not connected (unless $`𝒳`$ is a connected algebraic space). We write
$$I𝒳=\underset{i}{}𝒳_i$$
for the decomposition of $`I𝒳`$ into a disjoint union of connected components. Here $``$ is an index set. Among all components there is a distinguished one (indexed by $`0`$)
$$𝒳_0:=\{(x,id)|xOb(𝒳),idAut(x)\text{ is the identity element}\},$$
which is isomorphic to $`𝒳`$.
There is a natural involution $`I:I𝒳I𝒳`$ defined by interchanging the factors of $`𝒳\times _{𝒳\times 𝒳}𝒳`$. On objects we have $`I((x,g))=(x,g^1)`$. The restriction of $`I`$ to $`𝒳_i`$ is denoted by $`I_i`$. The map $`I_i`$ is an isomorphism between $`𝒳_i`$ and another component which we denote by $`𝒳_{i^I}`$. It is clear that $`𝒳_{i_{}^{I}{}_{}{}^{I}}=𝒳_i`$. Also, the restriction of $`I`$ to the distinguished component $`𝒳_0`$ is the identity map $`𝒳_0𝒳_0`$.
There is a locally constant function $`ord:I𝒳`$ defined by sending $`(x,g)`$ to the order of $`g`$ in $`Aut_𝒳(x)`$. Let $`r_i`$ denote its value on the connected component $`𝒳_i`$. Note that $`r_{i^I}=r_i`$. If we view $`I𝒳`$ as in Remark 2.1.2 (i), it is easy to see that the value of $`ord`$ at $`[B\mu _r𝒳]`$ is $`r`$.
###### Example 2.1.3.
Let $`𝒳`$ be of the form $`[M/G]`$ with $`M`$ a smooth variety and $`G`$ a finite group. We can take the index set $``$ to be the set $`\{(g)|gG\}`$ of conjugacy classes of $`G`$. In this case the centralizer $`C_G(g)`$ acts on the locus $`M^g`$ of $`g`$-fixed points. For the conjugacy class $`(g)`$ we have the component
$$𝒳_{(g)}=[M^g/C_G(g)],$$
and the distinguished component is $`[M^{id}/C_G(id)]=[M/G]`$. The morphism $`I_{(g)}`$ is an isomorphism between $`𝒳_{(g)}`$ and $`𝒳_{(g^1)}`$. In our notation, $`(g)^I=(g^1)`$. Also, the value of the function $`ord`$ on the component $`[M^g/C_G(g)]`$ is the order of the element $`g`$ in $`G`$.
### 2.2. Vector bundles on orbifolds
Let $`F`$ be a vector bundle on $`𝒳`$. When we view $`𝒳`$ as a geometric object locally a quotient of an affine scheme by a finite group, we may view $`F`$ as an object on $`𝒳`$ locally an equivariant vector bundle on an affine scheme. In this section we discuss some properties of the pullback bundle $`q^{}F`$, which is a vector bundle on the inertia stack $`I𝒳`$.
Denote by $`(q^{}F)_i`$ the restriction to $`𝒳_i`$ of the pullback of $`F`$, i.e. $`(q^{}F)_i:=q^{}F|_{𝒳_i}`$. At a point $`(x,g)𝒳_i`$, the fiber of $`(q^{}F)_i`$ admits an action of $`g`$, and is accordingly decomposed into a direct sum of eigenspaces of the $`g`$-action. This gives a global decomposition (see ),
$$(q^{}F)_i=\underset{0l<r_i}{}F_i^{(l)},$$
where $`F_i^{(l)}`$ is the eigen-subbundle with eigenvalue $`\zeta _{r_i}^l`$ and $`\zeta _{r_i}=\mathrm{exp}(2\pi \sqrt{1}/r_i)`$ is a primitive $`r_i`$-th root of unity. We make the convention that $`0l<r_i`$. Define $`(q^{}F)_i^{inv}:=F_i^{(0)}`$. Denote by $`q^{}F^{inv}`$ the bundle over $`I𝒳`$ whose restriction to $`𝒳_i`$ is $`(q^{}F)_i^{inv}`$.
The following result addresses compatibility of the decomposition of $`(q^{}F)_i`$ with pulling back by the involution $`I_i:𝒳_i𝒳_{i^I}`$.
###### Lemma 2.2.1.
1. $`I_i^{}(F_{i^I}^{(r_il)})=F_i^{(l)}`$ for $`0<l<r_i`$.
2. $`I_i^{}(F_{i^I}^{(0)})=F_i^{(0)}`$.
###### Proof.
We verify (1). Let $`S`$ be a scheme and $`(x,g^1)`$ an $`S`$-valued point of $`𝒳_{i^I}`$. Denote by $`\stackrel{~}{x}:S𝒳_{i^I}`$ the morphism corresponding to $`(x,g^1)`$. Then the $`S`$-valued point $`(x,g)`$ of $`𝒳_i`$ corresponds to the morphism $`I_{i^I}\stackrel{~}{x}:S𝒳_i`$.
Since $`F_{i^I}^{(r_il)}|_{(x,g^1)}:=\stackrel{~}{x}^{}(F_{i^I}^{(r_il)})`$ is the subbundle of $`F_{i^I}|_{(x,g^1)}:=\stackrel{~}{x}^{}F`$ on which $`g^1`$ acts with eigenvalue $`\zeta _{r_i}^{r_il}`$, $`g`$ acts on $`F_{i^I}^{(r_il)}|_{(x,g^1)}`$ with eigenvalue $`\zeta _{r_i}^l`$. Also,
$$(I_i^{}(F_{i^I}^{(r_il)}))|_{(x,g)}:=(I_{i^I}\stackrel{~}{x})^{}I_i^{}(F_{i^I}^{(r_il)})=\stackrel{~}{x}^{}(F_{i^I}^{(r_il)})=:F_{i^I}^{(r_il)}|_{(x,g^1)}.$$
Therefore $`(I_i^{}(F_{i^I}^{(r_il)}))|_{(x,g)}`$ is the subbuundle of $`\stackrel{~}{x}^{}F`$ on which $`g`$ acts with eigenvalue $`\zeta _{r_i}^l`$, which is $`F_i^{(l)}|_{(x,g)}:=(I_{i^I}\stackrel{~}{x})^{}(F_i^{(l)})`$. Hence $`I_i^{}(F_{i^I}^{(r_il)})F_i^{(l)}`$. The same argument proves that $`I_{i^I}^{}(F_i^{(l)})F_{i^I}^{(r_il)}`$. Since $`I_{i^I}I_i`$ is the identity map, we find $`F_i^{(l)}=I_i^{}I_{i^I}^{}(F_i^{(l)})I_i^{}(F_{i^I}^{(r_il)})`$. Thus $`I_i^{}(F_{i^I}^{(r_il)})=F_i^{(l)}`$.
A similar argument proves (2). ∎
We can describe the vector bundles $`F_i^{(l)}`$ using the identification (2.1.2.1). Each component $`𝒳_i`$ of $`I𝒳`$ can be viewed as the moduli stack of representable morphisms from constant $`\mu _{r_i}`$-gerbes to $`𝒳`$. Hence there is a universal family over $`𝒳_i`$:
$$\begin{array}{ccc}𝒳_i\times B\mu _{r_i}& \stackrel{\rho }{}& 𝒳\\ & & \\ 𝒳_i.\end{array}$$
Let $`\gamma :𝒳_i𝒳_i\times B\mu _{r_i}`$ be the morphism such that the map $`𝒳_i𝒳_i`$ to the first factor is the identity and the map $`𝒳_iB\mu _{r_i}`$ to the second factor<sup>7</sup><sup>7</sup>7In other words, the map $`𝒳_iB\mu _{r_i}`$ to the second factor is the composition $`𝒳_i\text{Spec}B\mu _{r_i}`$ where $`𝒳_i\text{Spec}`$ is the constant map and $`\text{Spec}B\mu _{r_i}`$ is the atlas of $`B\mu _{r_i}`$. corresponds to the trivial $`\mu _{r_i}`$-bundle over $`𝒳_i`$. The pull-back $`\rho ^{}F`$ admits an action of $`𝔲_{r_i}`$. Let $`(\rho ^{}F)^{(l)}`$ be the eigen sub-bundle of $`\rho ^{}F`$ on which $`𝔲_{r_i}`$ acts with eigenvalue $`\zeta _{r_i}^l`$. Then<sup>8</sup><sup>8</sup>8Note that the bundle $`(\rho ^{}F)^{(l)}`$ over $`𝒳_i\times B\mu _{r_i}`$ can be viewed as a bundle over $`𝒳_i`$ with a $`\mu _{r_i}`$-action. In this point of view pulling back by $`\gamma `$ simply forgets the $`\mu _{r_i}`$-action. we have
(2.2.1.1)
$$\gamma ^{}((\rho ^{}F)^{(l)})=F_i^{(l)}.$$
### 2.3. Orbifold Cohomology and Orbifold Cup Product
In this section we collect some facts about orbifold cohomology which we will use.
###### Definition 2.3.1.
Following Chen-Ruan , the cohomology $`H^{}(I𝒳,)`$ of the inertia stack is called the orbifold cohomology.
###### Remark 2.3.2.
In general, the cohomology with rational coefficients of a stack can be defined as the (singular) cohomology of a geometric realization of the simplicial scheme associated to this stack. For our purpose we define the cohomology of a Deligne-Mumford stack as the (singular) cohomology of its coarse moduli space. In our setting these two definitions are equivalent. See , Section 2.2 for a detailed discussion.
#### Grading on Orbifold Cohomology
According to (see also , ), the orbifold cohomology
$$H^{}(I𝒳,)=_iH^{}(𝒳_i,)$$
is equipped with a grading different from the usual one. This grading is explained below.
###### Definition 2.3.3.
For each component $`𝒳_i`$ of $`I𝒳`$, the age $`age(𝒳_i)`$ is defined as follows: Let $`(x,g)𝒳_i`$. The tangent space $`T_x𝒳`$ is decomposed into a direct sum $`_{0l<r_i}V^{(l)}`$ of eigenspaces according to the $`g`$-action, where $`V^{(l)}`$ is the eigenspace with eigenvalue $`\zeta _{r_i}^l`$, $`0l<r_i`$, and $`\zeta _{r_i}=\mathrm{exp}(2\pi \sqrt{1}\frac{1}{r_i})`$. The age is defined to be
$$age(𝒳_i):=\frac{1}{r_i}\underset{0l<r_i}{}ldim_{}V^{(l)}.$$
It is easy to see that this definition is independent of choices of $`(x,g)𝒳_i`$.
The following Lemma follows directly from the definition.
###### Lemma 2.3.4 (, Lemma 3.2.1).
$$age(𝒳_i)+age(𝒳_{i^I})=dim_{}𝒳dim_{}𝒳_i.$$
###### Definition 2.3.5.
The orbifold degree of a class $`aH^{}(𝒳_i,)`$ is defined to be
$$orbdeg(a):=deg(a)+2age(𝒳_i).$$
The orbifold degree gives a grading on $`H^{}(I𝒳,)`$ different from the usual one.
#### Orbifold Poincaré pairing
Following , Section 3.3, the orbifold Poincaré pairing
$$(,)_{orb}:H^{}(I𝒳,)\times H^{}(I𝒳,)$$
is defined as follows: For $`aH^{}(𝒳_i,),bH^{}(𝒳_{i^I},)`$, define
$$(a,b)_{orb}:=_{𝒳_i}aI_i^{}b,$$
where $`_{𝒳_i}`$ stands for the pushforward map $`H^{}(𝒳_i,)H^{}(\text{Spec },)`$. For other choices of classes $`a,b`$ supported on components of $`I𝒳`$ the pairing $`(a,b)_{orb}`$ is defined to be $`0`$. Obviously this definition extends linearly to a pairing on $`H^{}(I𝒳,)`$.
The orbifold Poincaré pairing pairs cohomology classes from a component $`𝒳_i`$ with classes from the isomorphic component $`𝒳_{i^I}`$. The fact that it is a non-degenerate pairing follows from the fact that the usual Poincaré pairing on $`H^{}(𝒳_i,)`$ is non-degenerate.
###### Definition 2.3.6.
In what follows we often fix a homogeneous additive basis $`\{\varphi _\alpha \}`$ of $`H^{}(I𝒳,)`$ such that each $`\varphi _\alpha `$ is supported on one component $`𝒳_i`$ of $`I𝒳`$. We denote by $`\{\varphi ^\alpha \}`$ the dual basis under orbifold Poincaré pairing.
#### Orbifold Cup Product
On $`H^{}(I𝒳,)`$ there is a product structure, defined in and , called the orbifold cup product (or Chen-Ruan cup product), which is different from the ordinary cup product on $`H^{}(I𝒳,)`$.
###### Definition 2.3.7.
For $`a,bH^{}(I𝒳,)`$, their orbifold cup product $`a_{orb}b`$ is defined as follows: For $`cH^{}(I𝒳,)`$,
$$(a_{orb}b,c)_{orb}:=a,b,c_{0,3,0},$$
where the right side is defined in Section 2.5.2.
Together with the grading by orbifold degrees, $`(H^{}(I𝒳,),_{orb})`$ is a graded $``$-algebra with unit $`\mathrm{𝟏}H^0(𝒳)`$.
In the following special case, we can compare the orbifold cup product with the ordinary cup product of $`H^{}(I𝒳,)`$.
###### Lemma 2.3.8.
For $`aH^{}(𝒳,)`$ and $`bH^{}(𝒳_i,)`$, the orbifold cup product $`a_{orb}b`$ is equal to the ordinary product $`q^{}ab`$ in $`H^{}(I𝒳,)`$.
###### Proof.
Using the identification
$$\overline{}_{0,3}(𝒳,0;i,i^I,0)𝒳_i\times B\mu _{r_i}\times B\mu _{r_i}$$
described in Remark 2.4.3 below, we find that $`a_{orb}bH^{}(𝒳_i,)`$. For $`cH^{}(𝒳_{i^I},)`$, by Definition 2.3.7 we have
$$(a_{orb}b,c)_{orb}=_{𝒳_i}q^{}abI_i^{}c.$$
On the other hand, by definition of the orbifold Poincaré pairing we have
$$(q^{}ab,c)_{orb}=_{𝒳_i}(q^{}ab)I_i^{}c.$$
We conclude by the non-degeneracy of the pairing $`(,)_{orb}`$. ∎
### 2.4. Moduli of Orbifold Stable Maps
In this section we discuss some properties of the moduli stacks of orbifold stable maps. We also set up notations used throughout the paper.
Let $`\overline{}_{g,n}(𝒳,d)`$ be the moduli stack of $`n`$-pointed genus $`g`$ orbifold stable maps to $`𝒳`$ of degree $`d`$ with sections to all gerbes (see , Section 4.5). The stack $`\overline{}_{g,n}(𝒳,d)`$ parametrizes the following objects:
$$\begin{array}{ccc}(𝒞,\{\mathrm{\Sigma }_i\})& \stackrel{𝔣}{}& 𝒳\\ & & \\ T,\end{array}$$
where
1. $`𝒞/T`$ is a prestable genus $`g`$ balanced nodal orbicurve<sup>9</sup><sup>9</sup>9In and , this is called a balanced twisted curve.,
2. for $`i=1,\mathrm{},n`$, the substack $`\mathrm{\Sigma }_i𝒞`$ is an étale cyclotomic gerbe over $`T`$ with a section (hence a trivialization), and
3. $`𝔣`$ is a representable morphism whose induced map between coarse moduli spaces is a $`n`$-pointed genus $`g`$ stable map of degree $`𝔣_{}[𝒞]=d\text{Eff}(𝒳)`$. (The object $`\text{Eff}(𝒳)`$ is defined in Definition 2.5.4 below. See , Section 2.2 for the definition of $`𝔣_{}`$.)
A precise definition of balanced nodal orbicurves can be found in and . The key idea is that an orbicurve is not a curve but a “stacky” version of curve: Nontrivial stack structures occur only at marked points or nodes. Étale locally near a marked point, an orbicurve over $`\text{Spec}`$ is isomorphic to the quotient $`[\text{Spec}[z]/\mu _r]`$ for some $`r`$, where the cyclic group $`\mu _r`$ acts on $`\text{Spec}[z]`$ via $`z\xi z,\xi \mu _r`$. Étale locally near a node, an orbicurve over $`\text{Spec}`$ is isomorphic to the quotient $`[\text{Spec}([x,y]/(xy))/\mu _r]`$ for some $`r`$, where $`\mu _r`$ acts on $`\text{Spec}([x,y]/(xy))`$ via $`x\xi x,y\xi ^1y,\xi \mu _r`$.
An étale cyclotomic gerbe over $`T`$ with a section is identified (via the trivialization given by the section) with $`T\times B\mu _r`$, where $`B\mu _r[\text{Spec}/\mu _r]`$ is the classifying stack associated to the finite group $`\mu _r`$.
###### Remark 2.4.1.
In and , orbifold stable maps are called twisted stable maps. Since the word “twisted” is used in a different context in this paper, we use the term “orbifold stable maps” instead.
For each $`i=1,\mathrm{},n`$ there is an evaluation map $`ev_i:\overline{}_{g,n}(𝒳,d)I𝒳`$ defined as follows: $`ev_i`$ sends an object $`𝔣:(𝒞,\{\mathrm{\Sigma }_i\})𝒳`$ to $`𝔣|_{\mathrm{\Sigma }_i}:\mathrm{\Sigma }_i𝒳`$. Since $`𝔣|_{\mathrm{\Sigma }_i}`$ is a map from a constant cyclotomic gerbe $`T\times B\mu _r`$ to $`𝒳`$, $`𝔣|_{\mathrm{\Sigma }_i}`$ is an object in $`I𝒳`$ by the description of $`I𝒳`$ in Remark 2.1.2 (i). We obtain a morphism $`ev_i:\overline{}_{g,n}(𝒳,d)I𝒳`$.
The stack $`\overline{}_{g,n}(𝒳,d)`$ can be decomposed according to images of the evaluation maps. Define
$$\overline{}_{g,n}(𝒳,d;i_1,\mathrm{},i_n):=ev_1^1(𝒳_{i_1})\mathrm{}ev_n^1(𝒳_{i_n})=_{j=1}^nev_j^1(𝒳_{i_j}).$$
We have
$$\overline{}_{g,n}(𝒳,d)=\underset{i_1,\mathrm{},i_n}{}\overline{}_{g,n}(𝒳,d;i_1,\mathrm{},i_n).$$
This decomposition according to images of the evaluation maps will be important for us: later on in our computations we will need explicit control on stack structures at the marked points.
The universal family over the moduli stack $`\overline{}_{g,n}(𝒳,d)`$ also admits a modular description. Let
$$\overline{}_{g,n+1}(𝒳,d)^{}:=ev_{n+1}^1(𝒳_0)\overline{}_{g,n+1}(𝒳,d)$$
denote the open-and-closed substack consisting of orbifold stable maps with trivial stack structure on the $`(n+1)`$-st marked point. According to , Corollary 9.1.3, there is a morphism
$$f:\overline{}_{g,n+1}(𝒳,d)^{}\overline{}_{g,n}(𝒳,d)$$
which forgets the $`(n+1)`$-st marked point. Moreover, $`f`$ exhibits $`\overline{}_{g,n+1}(𝒳,d)^{}`$ as the universal family over $`\overline{}_{g,n}(𝒳,d)`$, and $`ev_{n+1}:\overline{}_{g,n+1}(𝒳,d)^{}𝒳_0𝒳`$ is the universal orbifold stable map. Similarly, the universal family over the substack $`\overline{}_{g,n}(𝒳,d;i_1,\mathrm{},i_n)`$ is
$$\overline{}_{g,n+1}(𝒳,d;i_1,\mathrm{},i_n,0)\overline{}_{g,n}(𝒳,d;i_1,\mathrm{},i_n).$$
###### Remark 2.4.2.
There is another moduli stack $`𝒦_{g,n}(𝒳,d)`$ studied in , which parametrizes orbifold stable maps without trivializing the gerbes. Over $`𝒦_{g,n}(𝒳,d)`$ there are $`n`$ universal gerbes $`𝔖_j,1jn`$, corresponding to the marked points, and the fiber product over $`𝒦_{g,n}(𝒳,d)`$ of these gerbes is $`\overline{}_{g,n}(𝒳,d)`$, see , Section 4.5. We will not use this stack $`𝒦_{g,n}(𝒳,d)`$ to construct orbifold Gromov-Witten theory.
###### Remark 2.4.3.
We discuss briefly the special case $`(g,n,d)=(0,3,0)`$. According to , Section 6.2, the evaluation maps of $`𝒦_{0,3}(𝒳,0)`$ can be taken to have target $`I𝒳`$. Moreover, by , Lemma 7.7, $`𝒦_{0,3}(𝒳,0)^{}`$ is isomorphic to $`I𝒳`$. See also the proof of , Proposition 8.2.1. Under this isomorphism, $`ev_1`$ is identified with the identity map $`I𝒳I𝒳`$, $`ev_2`$ is identified with $`I:I𝒳I𝒳`$, and $`ev_3`$ is identified with $`q:I𝒳𝒳`$. The space $`\overline{}_{0,3}(𝒳,0;i_1,i_2,0)`$ is empty if $`i_2i_1^I`$. We have an isomorphism
$$\overline{}_{0,3}(𝒳,0;i,i^I,0)𝒳_i\times B\mu _{r_i}\times B\mu _{r_i}.$$
#### 2.4.4. Marked Points and Nodes
The marked points define divisors in the universal family $`\overline{}_{g,n+1}(𝒳,d)^{}`$. Let
$$𝒟_j\overline{}_{g,n+1}(𝒳,d)^{}$$
be the $`j`$-th universal gerbe over $`\overline{}_{g,n}(𝒳,d)`$. By definition, $`𝒟_j`$ is the pullback to $`\overline{}_{g,n}(𝒳,d)`$ of the gerbe $`𝔖_j`$ over $`𝒦_{g,n}(𝒳,d)`$. Since $`\overline{}_{g,n}(𝒳,d)`$ is the fiber product of all the $`𝔖_j`$’s, the pullback gerbe $`𝒟_j`$ admits a canonical section and is thus trivialized by this section. So for each $`1jn`$ there is a section $`\overline{}_{g,n}(𝒳,d)\overline{}_{g,n+1}(𝒳,d)^{}`$ corresponding to the $`j`$-th marked point. The image of this section is $`𝒟_j`$.
The identification of $`\overline{}_{g,n+1}(𝒳,d)^{}`$ as the universal family over $`\overline{}_{g,n}(𝒳,d)`$ implies that $`𝒟_j`$ can be described as a moduli space parametrizing maps $`𝔣:(𝒞,\{\mathrm{\Sigma }_i\})𝒳`$ with the following property: the domain has a distinguished balanced node $`\mathrm{\Sigma }𝒞`$ separating two parts $`𝒞_0`$ and $`𝒞_1`$. The marked points $`\mathrm{\Sigma }_j`$ and $`\mathrm{\Sigma }_{n+1}`$ lie on $`𝒞_1`$ and the other marked points lie on $`𝒞_0`$. $`𝔣|_{𝒞_0}:(𝒞_0,\{\mathrm{\Sigma }_i\}_{ij,n+1},\mathrm{\Sigma })𝒳`$ is an $`n`$-pointed orbifold stable map of genus $`g`$ and degree $`d`$, and $`𝔣|_{𝒞_1}:(𝒞_1,\mathrm{\Sigma },\mathrm{\Sigma }_j,\mathrm{\Sigma }_{n+1})𝒳`$ is a $`3`$-pointed orbifold stable map of genus $`0`$ and degree $`0`$.
Put
$$𝒟_{j,(i_1,\mathrm{},i_n)}:=𝒟_j\overline{}_{g,n+1}(𝒳,d;i_1,\mathrm{},i_n,0).$$
$`𝒟_{j,(i_1,\mathrm{},i_n)}`$ is the $`j`$-th universal gerbe over $`\overline{}_{g,n}(𝒳,d;i_1,\mathrm{},i_n)`$, which according to the discussion above is canonically trivialized. We have that $`𝒟_{j,(i_1,\mathrm{},i_n)}`$ is isomorphic to $`\overline{}_{g,n}(𝒳,d;i_1,\mathrm{},i_n)\times B\mu _{r_{i_j}}`$. Under this isomorphism, $`f|_{𝒟_{j,(i_1,\mathrm{},i_n)}}`$ coincides with the projection to the first factor.
Let $`𝒵\overline{}_{g,n+1}(𝒳,d)^{}`$ be the locus of nodes in the universal family. $`𝒵`$ is a disjoint union
$$𝒵=𝒵^{irr}𝒵^{red},$$
where $`𝒵^{irr}`$ is the locus of non-separating nodes and $`𝒵^{red}`$ is the locus of separating nodes. $`𝒵`$ is of virtual codimension two in $`\overline{}_{g,n+1}(𝒳,d)^{}`$, and is a cyclotomic gerbe over $`f(𝒵)`$.
There is a locally constant function $`ord:𝒵`$ defined by assigning to a node the order of its automorphism group: If a node is locally the quotient $`[U/\mu _r]`$ where $`U`$ is the curve $`xy=t`$ and the cyclic group $`\mu _r`$ of order $`r`$ acts via $`(x,y)(\zeta x,\zeta ^1y)`$, $`\zeta \mu _r`$, then $`ord`$ sends this node to the integer $`r`$.
Let $`𝒵_r:=ord^1(r)𝒵`$. Define
$$𝒵_{(i_1,\mathrm{},i_n)}:=𝒵\overline{}_{g,n+1}(𝒳,d;i_1,\mathrm{},i_n,0),𝒵_{r,(i_1,\mathrm{},i_n)}:=𝒵_r\overline{}_{g,n+1}(𝒳,d;i_1,\mathrm{},i_n,0).$$
The substacks $`𝒵_{(i_1,\mathrm{},i_n)}^{irr},𝒵_{(i_1,\mathrm{},i_n)}^{red},𝒵_{r,(i_1,\mathrm{},i_n)}^{irr},𝒵_{r,(i_1,\mathrm{},i_n)}^{red}\overline{}_{g,n+1}(𝒳,d;i_1,\mathrm{},i_n,0)`$ are similarly defined.
#### 2.4.5. Stable maps to the coarse moduli space
Let $`\overline{}_{g,n}(X,d)`$ be the moduli stack of $`n`$-pointed genus $`g`$ stable maps of degree $`d`$ to the coarse moduli space $`X`$. The universal family over $`\overline{}_{g,n}(X,d)`$ is $`\overline{}_{g,n+1}(X,d)\overline{}_{g,n}(X,d)`$, see for example , Corollary 4.6. There is a morphism
$$\pi _n:\overline{}_{g,n}(𝒳,d)\overline{}_{g,n}(X,d),$$
which sends an orbifold stable map to its induced stable map between coarse moduli spaces, see , Theorem 1.4.1.
### 2.5. Orbifold Gromov-Witten Theory
In this section we describe the Gromov-Witten theory for Deligne-Mumford stacks following , which is based on the stacks $`\overline{}_{g,n}(𝒳,d)`$. We refer the reader to and for a construction of orbifold Gromov-Witten theory based on the stacks $`𝒦_{g,n}(𝒳,d)`$ (see Remark 2.4.2). Both constructions yield the same orbifold Gromov-Witten invariants. The intersection theory for algebraic stacks needed here can be found in and (which has already been used to construct Gromov-Witten theory for varieties).
#### 2.5.1. Virtual fundamental classes and descendants
The stack $`\overline{}_{g,n}(𝒳,d)`$ admits a perfect obstruction theory relative to the Artin stack of prestable pointed orbicurves (, Section 4.6). This obstruction theory is given by the object $`(𝐑^{}f_{}ev_{n+1}^{}T𝒳)^{}`$ in the derived category $`D(Coh(\overline{}_{g,n}(𝒳,d)))`$. Results in and apply to yield a virtual fundamental class
$$[\overline{}_{g,n}(𝒳,d)]^{vir}H_{}(\overline{}_{g,n}(𝒳,d),).$$
The virtual fundamental class $`[\overline{}_{g,n}(𝒳,d;i_1,\mathrm{},i_n)]^{vir}`$ may be obtained by restriction. As observed in , we need to use a weighted virtual fundamental class $`[\overline{}_{g,n}(𝒳,d)]^w`$ defined as follows: the restriction of $`[\overline{}_{g,n}(𝒳,d)]^w`$ to $`\overline{}_{g,n}(𝒳,d;i_1,\mathrm{},i_n)`$, which we denote by $`[\overline{}_{g,n}(𝒳,d;i_1,\mathrm{},i_n)]^w`$, is defined by
$$[\overline{}_{g,n}(𝒳,d;i_1,\mathrm{},i_n)]^w:=(\underset{j=1}{\overset{n}{}}r_{i_j})[\overline{}_{g,n}(𝒳,d;i_1,\mathrm{},i_n)]^{vir}.$$
We refer to , Section 4.6 for more details.
We now define the descendant classes. For each $`i=1,\mathrm{},n`$, the universal family $`\overline{}_{g,n+1}(X,d)\overline{}_{g,n}(X,d)`$ has a section
$$\sigma _i:\overline{}_{g,n}(X,d)\overline{}_{g,n+1}(X,d),$$
which corresponds to the $`i`$-th marked point (note that here we consider the moduli stack of stable maps to the coarse moduli space $`X`$). Recall that the $`i`$-th tautological line bundle is defined to be the pullback of the relative dualizing sheaf $`\omega _{\overline{}_{g,n+1}(X,d)/\overline{}_{g,n}(X,d)}`$ by $`\sigma _i`$,
$$L_i:=\sigma _i^{}\omega _{\overline{}_{g,n+1}(X,d)/\overline{}_{g,n}(X,d)},$$
see for example . Let $`\psi _i=c_1(L_i)`$ and
$$\overline{\psi }_i:=\pi _n^{}\psi _iH^2(\overline{}_{g,n}(𝒳,d),).$$
These $`\overline{\psi }_i`$ are the descendant classes of $`\overline{}_{g,n}(𝒳,d)`$. Note that our choice of descendant classes differs from those of by constants.
#### 2.5.2. Untwisted Theory
We are now ready to define the invariants, following and . Let $`a_jH^{p_j}(𝒳_{i_j},),j=1,\mathrm{},n`$ be cohomology classes and $`k_1,\mathrm{},k_n`$ nonnegative integers. We define the descendant orbifold Gromov-Witten invariants to be
$$a_1\overline{\psi }^{k_1},\mathrm{},a_n\overline{\psi }^{k_n}_{g,n,d}:=_{[\overline{}_{g,n}(𝒳,d;i_1,\mathrm{},i_n)]^w}(ev_i^{}a_1)\overline{\psi }_1^{k_1}\mathrm{}(ev_n^{}a_n)\overline{\psi }_n^{k_n}.$$
Here $`_{[\overline{}_{g,n}(𝒳,d;i_1,\mathrm{},i_n)]^w}`$ stands for capping with the virtual fundamental class $`[\overline{}_{g,n}(𝒳,d;i_1,\mathrm{},i_n)]^w`$ followed by pushing forward to $`\text{Spec }`$. The symbol $`\mathrm{}_{g,n,d}`$ is by definition multi-linear in its entries.
The invariant $`a_1\overline{\psi }^{k_1},\mathrm{},a_n\overline{\psi }^{k_n}_{g,n,d}`$ is zero unless
(2.5.2.1)
$$\frac{1}{2}(orbdeg(a_1)+\mathrm{}+orbdeg(a_n))+k_1+\mathrm{}+k_n=(1g)(dim_{}𝒳3)+n+_dc_1(T_𝒳),$$
where $`orbdeg(a_j)=p_j+2age(𝒳_{i_j})`$ is the orbifold degree defined in Section 2.3. This follows from the formula for virtual dimension of $`\overline{}_{g,n}(𝒳,d;i_1,\mathrm{},i_n)`$, which follows from , Theorem 7.2.1.
###### Remark 2.5.3.
The cohomology $`H^{}(I𝒳,)`$ is viewed as a super vector space. For simplicity we systematically ignore the signs that may come out. It is straightforward to include the signs in our results (c.f. ).
We can form generating functions to encode these invariants.
###### Definition 2.5.4.
Let $`𝐭=𝐭(z)=t_0+t_1z+t_2z^2+\mathrm{}H^{}(I𝒳)[z]`$. Define
$$𝐭,\mathrm{},𝐭_{g,n,d}=𝐭(\overline{\psi }),\mathrm{},𝐭(\overline{\psi })_{g,n,d}:=\underset{k_1,\mathrm{},k_n0}{}t_{k_1}\overline{\psi }^{k_1},\mathrm{},t_{k_n}\overline{\psi }^{k_n}_{g,n,d}.$$
The total descendant potential is defined to be
$$𝒟_𝒳(𝐭):=\mathrm{exp}\left(\underset{g0}{}\mathrm{}^{g1}_𝒳^g(𝐭)\right),$$
where
$$_𝒳^g(𝐭):=\underset{n0,d\text{Eff}(𝒳)}{}\frac{Q^d}{n!}𝐭,\mathrm{},𝐭_{g,n,d}.$$
Here $`\mathrm{}`$ is a formal variable, and $`Q^d`$ is an element of the Novikov ring $`\mathrm{\Lambda }_{nov}`$ which is a completion of the group ring $`[\text{Eff}(𝒳)]`$ of the semi-group $`\text{Eff}(𝒳)`$ of effective curve classes (i.e. classes in $`H_2(𝒳,)`$ represented by images of representable maps from complete stacky curves to $`𝒳`$). The completion is done with respect to an additive valuation
$$v\left(\underset{d\text{Eff}(𝒳)}{}c_dQ^d\right)=\text{min}_{c_d0}_dc_1(L)$$
defined by the ample polarization $`L`$ of the coarse moduli space $`X`$ which we choose once and for all.
$`_𝒳^g(𝐭)`$ is called the genus-$`g`$ descendant potential. It is regarded as a $`\mathrm{\Lambda }_{nov}`$-valued formal power series in the variables $`t_k^\alpha `$ where
$$t_k=\underset{\alpha }{}t_k^\alpha \varphi _\alpha H^{}(I𝒳,),k0.$$
###### Remark 2.5.5.
Following the treatement in , Section 2.2, the homology group $`H_2(𝒳,)`$ with rational coefficients is defined to be the homology group $`H_2(X,)`$ of the coarse moduli space $`X`$. For this reason degree of effective curve classes in $`𝒳`$ are identified with degrees of effective curve classes in $`X`$, and we will use the term interchangably.
###### Lemma 2.5.6 (c.f. , Lemma 1.3.1).
$`𝒟_𝒳`$ is well-defined as a formal power series in the variables $`t_k^\alpha `$ taking values in $`\mathrm{\Lambda }_{nov}[[\mathrm{},\mathrm{}^1]]`$.
###### Proof.
We follows the argument of , Lemma 1.3.1, which treats the manifold case. First of all, the expression
(2.5.6.1)
$$\underset{g0}{}\mathrm{}^{g1}_𝒳^g(𝐭)$$
is well-defined as a formal power series in $`t_k^\alpha `$ taking values in $`\mathrm{\Lambda }_{nov}[[\mathrm{},\mathrm{}^1]]`$. Given a monomial $`\mathrm{}^{g1}Q^d_{1in}(t_{k_i}^{\alpha _i})^{j_i}`$, we define its degree to be the triple $`(g1,_{1in}j_i,d)`$. The coefficient of a monomial of degree $`(a,b,c)`$ that occurs in (2.5.6.1) is a (non-zero) orbifold Gromov-Witten invariant coming from the moduli space $`\overline{}_{a+1,b}(𝒳,c)`$. One observes that
1. Since $`\overline{}_{a+1,b}(𝒳,c)`$ is finite dimensional, in each degree only finitely many monomials can occur in (2.5.6.1);
2. Since $`\overline{}_{0,0}(𝒳,0)`$ and $`\overline{}_{1,0}(𝒳,0)`$ are empty, if a monomial of degree $`(a,b,0)`$ occurs in (2.5.6.1), then at least one of $`a`$ and $`b`$ is strictly positive.
Now, a monomial of degree $`(a,b,c)`$ occurs in $`𝒟_𝒳=\mathrm{exp}(_{g0}\mathrm{}^{g1}_𝒳^g(𝐭))`$ only if there are monomials of degrees $`(a_1,b_1,c_1),\mathrm{},(a_N,b_N,c_N)`$ in (2.5.6.1) such that
$$a_1+\mathrm{}+a_N=a,b_1+\mathrm{}+b_N=b,c_1+\mathrm{}+c_N=c.$$
By the observations above, there are only finitely many choices of $`\{(a_i,b_i,c_i)\}`$. The result follows. ∎
We remark that the orbifold Gromov-Witten theory considered here differs slightly from those in and : we work with algebraic stacks while works with symplectic orbifolds. But unlike , we work with cohomology instead of Chow ring. One reason for this is that Poincaré duality holds for cohomology, but not for Chow rings in general. A definition of cohomological orbifold Gromov-Witten invariants of Deligne-Mumford stacks using the moduli stack $`𝒦_{g,n}(𝒳,d)`$ can be found in . This definition is equivalent to ours.
#### 2.5.7. Universal equations in orbifold Gromov-Witten theory
The Gromov-Witten invariants for varieties are known to satisfy four sets of universal equations<sup>10</sup><sup>10</sup>10 These universal equations can be rewritten as differential equations of the generating functions.: string equation (SE), divisor equation (DIV), dilaton equation (DE), and topological recursion relations (TRR). One may find the precise forms of these equations in for instance . The proof of these equations is based on comparisons of descendant classes on various moduli spaces related by forgetful maps. These four sets of equations hold in orbifold Gromov-Witten theory as well, and they take the same form as those in Gromov-Witten theory for varieties. More precisely, we have
String equation:
$$a_1\overline{\psi }^{k_1},\mathrm{},a_n\overline{\psi }^{k_n},\mathrm{𝟏}_{g,n+1,d}=\underset{j=1}{\overset{n}{}}a_1\overline{\psi }^{k_1},\mathrm{},a_j\overline{\psi }^{k_j1},\mathrm{},a_n\overline{\psi }^{k_n}_{g,n,d};$$
Divisor equation: for $`\gamma H^2(𝒳,)`$,
$$\begin{array}{cc}& a_1\overline{\psi }^{k_1},\mathrm{},a_n\overline{\psi }^{k_n},\gamma _{g,n+1,d}=\left(_d\gamma \right)a_1\overline{\psi }^{k_1},\mathrm{},a_n\overline{\psi }^{k_n}_{g,n,d}\hfill \\ & +\underset{j=1}{\overset{n}{}}a_1\overline{\psi }^{k_1},\mathrm{},(\gamma _{orb}a_j)\overline{\psi }^{k_j1},\mathrm{},a_n\overline{\psi }^{k_n}_{g,n,d};\hfill \end{array}$$
Dilaton equation:
$$a_1\overline{\psi }^{k_1},\mathrm{},a_n\overline{\psi }^{k_n},\mathrm{𝟏}\overline{\psi }_{g,n+1,d}=(2g2+n)a_1\overline{\psi }^{k_1},\mathrm{},a_n\overline{\psi }^{k_n}_{g,n,d};$$
Topological recursion relations (in genus zero): for $`tH^{}(I𝒳)`$, define
$$a_1\overline{\psi }^{k_1},\mathrm{},a_n\overline{\psi }^{k_n}_0:=\underset{k=0}{\overset{\mathrm{}}{}}\underset{d\text{Eff}(𝒳)}{}\frac{Q^d}{k!}a_1\overline{\psi }^{k_1},\mathrm{},a_n\overline{\psi }^{k_n},t,\mathrm{},t_{0,n+k,d}.$$
Then
$$\varphi _{\alpha _1}\overline{\psi }^{k_1+1},\varphi _{\alpha _2}\overline{\psi }^{k_2},\varphi _{\alpha _3}\overline{\psi }^{k_3}_0=\underset{\alpha }{}\varphi _{\alpha _1}\overline{\psi }^{k_1},\varphi _\alpha _0\varphi ^\alpha ,\varphi _{\alpha _2}\overline{\psi }^{k_2},\varphi _{\alpha _3}\overline{\psi }^{k_3}_0,$$
where $`\{\varphi _\alpha \}`$ is an additive basis of $`H^{}(I𝒳)`$ and $`\{\varphi ^\alpha \}`$ its dual basis under orbifold Poincaré pairing. In these equations the term $`\overline{\psi }^1`$ is defined to be $`0`$.
Proofs of (SE), (DIV) and (DE) can be found in . The key observation is that, since our choice of descendant classes are pulled back from moduli space of stable maps to the coarse moduli space, the comparisons of various descendant classes remain unchanged. See for more details. A proof of (TRR) is given in Appendix E.
#### 2.5.8. Twisted Theory
We now introduce twisted orbifold Gromov-Witten invariants. We will make the following
###### Assumption 2.5.9.
$`𝒳`$ is a quotient of a smooth quasi-projective scheme by a linear algebraic group.
Given a vector bundle $`F`$ over $`𝒳`$ and an invertible multiplicative characteristic class $`𝐜()=\mathrm{exp}(_ks_kch_k())`$. We define the “twisting factor” as follows.
###### Definition 2.5.10.
For a vector bundle $`F`$ on $`𝒳`$, define
$$F_{g,n,d}:=f_{}ev_{n+1}^{}F,$$
where $`f_{}`$ is the K-theoretic pushforward. Assumption 2.5.9 and the results of imply that the map $`f`$ is a local complete intersection morphism. Therefore the K-theoretic push-forward $`f_{}`$ of a bundle has a locally free resolution and thus defines an element in the Grothendieck group $`K^0`$. Hence $`F_{g,n,d}`$ is an element in $`K^0(\overline{}_{g,n}(𝒳,d))`$. Its restriction to $`\overline{}_{g,n}(𝒳,d;i_1,\mathrm{},i_n)`$, which is an element in $`K^0(\overline{}_{g,n}(𝒳,d;i_1,\mathrm{},i_n))`$, is denoted by $`F_{g,n,d,(i_1,..,i_n)}`$. The cohomology classes $`𝐜(F_{g,n,d})`$ and $`𝐜(F_{g,n,d,(i_1,..,i_n)})`$ are called the twisting factors.
More detailed discussions and properties of $`F_{g,n,d}`$ can be found in Appendix B.
We define the $`(𝐜,F)`$-twisted descendant orbifold Gromov-Witten invariants to be
$$a_1\overline{\psi }^{k_1},\mathrm{},a_n\overline{\psi }^{k_n}_{g,n,d,(𝐜,F)}:=_{[\overline{}_{g,n}(𝒳,d;i_1,\mathrm{},i_n)]^w}(ev_i^{}a_1)\overline{\psi }_1^{k_1}\mathrm{}(ev_n^{}a_n)\overline{\psi }_n^{k_n}𝐜(F_{g,n,d,(i_1,..,i_n)}),$$
where $`a_1,\mathrm{},a_n`$ are as in Section 2.5.2. The symbol $`\mathrm{}_{g,n,d,(𝐜,F)}`$ is by definition multi-linear in its entries. Again, these invariants can be packaged into generating functions.
###### Definition 2.5.11.
Define
$$𝐭,\mathrm{},𝐭_{g,n,d,(𝐜,F)}=𝐭(\overline{\psi }),\mathrm{},𝐭(\overline{\psi })_{g,n,d,(𝐜,F)}:=\underset{k_1,\mathrm{},k_n0}{}t_{k_1}\overline{\psi }^{k_1},\mathrm{},t_{k_n}\overline{\psi }^{k_n}_{g,n,d,(𝐜,F)}.$$
The $`(𝐜,F)`$-twisted total descendant potential is defined to be
$$𝒟_{(𝐜,F)}(𝐭):=\mathrm{exp}\left(\underset{g0}{}\mathrm{}^{g1}_{(𝐜,F)}^g(𝐭)\right),$$
where
$$_{(𝐜,F)}^g(𝐭):=\underset{n0,d\text{Eff}(𝒳)}{}\frac{Q^d}{n!}𝐭,\mathrm{},𝐭_{g,n,d,(𝐜,F)}.$$
$`_{(𝐜,F)}^g(𝐭)`$ is regarded as a formal power series in the variables $`t_k^\alpha `$ taking values in the ring $`\mathrm{\Lambda }_s`$. The ring $`\mathrm{\Lambda }_s`$ is defined to be the completion of $`[\text{Eff}(𝒳)][s_0,s_1,\mathrm{}]`$ with respect to the additive valuation
$$v\left(\underset{d\text{Eff}(𝒳)}{}c_dQ^d\right)=\text{min}_{c_d0}_dc_1(L),v(s_k)=k+1(\text{Here }L\text{ is the chosen ample line bundle on }X).$$
The total descendant potential $`𝒟_{(𝐜,F)}`$ is well-defined as a formal power series in $`t_k^\alpha `$ taking values in $`\mathrm{\Lambda }_s[[\mathrm{},\mathrm{}^1]]`$. The proof of Lemma 2.5.6 can be easily adjusted to treat this case.
## 3. Givental’s symplectic space formalism
A. Givental introduces a symplectic vector space formalism to describe Gromov-Witten theory (see , , ). In this formalism many properties of Gromov-Witten invariants can be studied using linear symplectic transformations of a certain symplectic vector space, making them more geometric. In this section we explain how this formalism is applied to orbifold Gromov-Witten theory. We will present this in detail for both twisted and untwisted theories.
To take care of certain convergence issues, we will make use of the following definition.
###### Definition 3.0.1 (c.f. ).
Let $`R`$ be a topological ring with an additive valuation $`v:R\{0\}`$. Define the space of $`R`$-valued convergent Laurent series in $`z`$ to be
$$R\{z,z^1\}:=\left\{\underset{n}{}r_nz^n\text{ : }r_nR,v(r_n)\mathrm{}\text{ as }|n|\mathrm{}\right\}.$$
Note that $`R\{z,z^1\}`$ is a ring if $`R`$ is complete. Also put
$$\begin{array}{c}\hfill R\{z\}:=\left\{\underset{n0}{}r_nz^n\text{ : }r_nR,v(r_n)\mathrm{}\text{ as }n\mathrm{}\right\},\\ \hfill R\{z^1\}:=\left\{\underset{n0}{}r_nz^n\text{ : }r_nR,v(r_n)\mathrm{}\text{ as }n\mathrm{}\right\}.\end{array}$$
### 3.1. Givental’s formalism for untwisted theory
Consider the space
$$:=H^{}(I𝒳,)\mathrm{\Lambda }_s\{z,z^1\}$$
of orbifold-cohomology-valued convergent Laurent series. There is a $`\mathrm{\Lambda }_s`$-valued symplectic form on $``$ given by
$$\mathrm{\Omega }(f,g)=\text{Res}_{z=0}(f(z),g(z))_{orb}dz,\text{for }f,g.$$
Consider the following polarization
(3.1.0.1)
$$\begin{array}{cc}& =_+_{},\hfill \\ & _+:=H^{}(I𝒳,)\mathrm{\Lambda }_s\{z\}\text{ and }_{}:=z^1H^{}(I𝒳,)\mathrm{\Lambda }_s\{z^1\}.\hfill \end{array}$$
This identifies $``$ with $`_+_+^{}`$, where $`_+^{}`$ is the dual $`\mathrm{\Lambda }_s`$-module. (We may thus think of $``$ as the cotangent bundle $`T^{}_+`$.) Both $`_+`$ and $`_{}`$ are Lagrangian subspaces with respect to $`\mathrm{\Omega }`$.
Introduce a Darboux coordinate system $`\{p_a^\mu ,q_b^\nu \}`$ on $`(,\mathrm{\Omega })`$ with respect to the polarization (3.1.0.1). Namely, in these coordinates, a general point in $``$ takes the form
$$\underset{a0}{}\underset{\mu }{}p_a^\mu \varphi ^\mu (z)^{a1}+\underset{b0}{}\underset{\nu }{}q_b^\nu \varphi _\nu z^b.$$
Put $`p_a=_\mu p_a^\mu \varphi ^\mu `$ and $`q_b=_\nu q_b^\nu \varphi _\nu `$. Denote
$$\begin{array}{cc}& 𝐩=𝐩(z):=\underset{k0}{}p_k(z)^{k1}=p_0(z)^1+p_1(z)^2+\mathrm{};\hfill \\ & 𝐪=𝐪(z):=\underset{k0}{}q_kz^k=q_0+q_1z+q_2z^2+\mathrm{}.\hfill \end{array}$$
For $`𝐭(z)_+=H^{}(I𝒳,)\mathrm{\Lambda }_s\{z\}`$ introduce a shift $`𝐪(z)=𝐭(z)\mathrm{𝟏}z`$ called the dilaton shift.
Define the Fock space Fock to be the space of formal functions<sup>11</sup><sup>11</sup>11This means formal power series in variables $`t_k^\alpha `$ where $`t_k=_\alpha t_k^\alpha \varphi _\alpha `$. in $`𝐭(z)_+`$ taking values in $`\mathrm{\Lambda }_s[[\mathrm{},\mathrm{}^1]]`$. In other words, Fock is the space of formal functions on $`_+`$ in the formal neighborhood of $`𝐪=\mathrm{𝟏}z`$. The descendant potential $`𝒟_𝒳(𝐭)`$ is regarded as an element in Fock via the dilaton shift.
The generating function $`_𝒳^0`$ of genus-$`0`$ orbifold Gromov-Witten invariants, which is defined in a formal neighborhood of $`\mathrm{𝟏}z`$, defines a formal germ of Lagrangian submanifold
$$_𝒳:=\{(𝐩,𝐪)|𝐩=d_𝐪_𝒳^0\}=T^{}_+,$$
which is just the graph of the differential of $`_𝒳^0`$. Equivalently $`_𝒳`$ is defined by all equations of the form $`p_a^\mu =\frac{_𝒳^0}{q_a^\mu }`$.
By , Theorem 1, string and dilaton equations and topological recursion relations imply that $`_𝒳`$ satisfies the following properties.
###### Theorem 3.1.1 (c.f. ).
$`_𝒳`$ is the formal germ of a Lagrangian cone with vertex at the origin such that each tangent space $`T`$ to the cone is tangent to the cone exactly along $`zT`$. In other words, if $`N`$ is a formal neighbourhood in $``$ of the unique geometric point on $`_𝒳`$, then
(3.1.1.1)
$$\begin{array}{cc}& \text{(1) }T_𝒳=zTN;\hfill \\ & \text{(2) }\text{for each }𝐟zTN,\text{ the tangent space to }_𝒳\text{ at }𝐟\text{ is }T;\hfill \\ & \text{(3) }\text{if }T=T_𝐟_𝒳\text{ then }fzTN.\hfill \end{array}$$
The statements in (3.1.1.1) are valid in the context of formal geometry. So for instance $`T_𝒳=zTN`$ means that any formal family of elements of $``$ which is both a family of elements of $`T`$ and of $`_𝒳`$ is also a family of elements of both $`zT`$ and $`N`$, and vice versa. Also, these statements imply that the tangent spaces $`T`$ of $`_𝒳`$ are closed under multiplication by $`z`$. Moreover, because $`T/zT`$ is isomorphic to $`H^{}(I𝒳,)`$, it follows from (3.1.1.1) that $`_𝒳`$ is the union of the (finite-dimensional) family of germs of (infinite-dimensional) linear subspaces
$$\{zTN|T\text{ is a tangent space of }_𝒳\}.$$
###### Definition 3.1.2.
Following , we define the $`J`$-function $`J_𝒳(t,z)`$ as follows,
$$J_𝒳(t,z)=z+t+\underset{n1,d\text{Eff}(𝒳)}{}\frac{Q^d}{(n1)!}\underset{k0,\alpha }{}t,\mathrm{},t,\varphi _\alpha \overline{\psi }^k_{0,n,d}\frac{\varphi ^\alpha }{z^{k+1}}.$$
This is a formal power series in coordinates $`t^\alpha `$ of $`t=_\alpha t^\alpha \varphi _\alpha H^{}(I𝒳,)`$ taking values in $``$. The point of $`_𝒳`$ above $`z+t_+`$ is $`J_𝒳(t,z)`$.
For each $`k0`$, the coefficient of the $`z^{1k}`$ term in $`J_𝒳(t,z)`$ takes values in $`H^{}(I𝒳,)\mathrm{\Lambda }_s`$. According to the decomposition $`H^{}(I𝒳,)=_iH^{}(𝒳_i,)`$, we write
$$J_𝒳(t,z)=(J_i(t,z))\text{where}J_i(t,z)\text{ takes values in }H^{}(𝒳_i,)\mathrm{\Lambda }_s\{z,z^1\}.$$
We further decompose $`J_i`$ according to degrees,
$$J_i(t,z)=\underset{d\text{Eff}(𝒳)}{}J_{i,d}(t,z)Q^d.$$
This $`J`$-function plays an important role in the genus-0 theory. For example:
###### Lemma 3.1.3.
The union of the (finite-dimensional) family
$$tzT_{J_𝒳(t,z)}_𝒳N,t\text{ in a formal neighborhood of zero in }H^{}(I𝒳,)\mathrm{\Lambda }_s,$$
of germs of linear subspaces is $`_𝒳`$.
###### Proof.
According to the discussion above, we just need to prove that every tangent space $`T`$ of $`_𝒳`$ is of the form $`T_{J_𝒳(\tau ,z)}_𝒳`$ for some $`\tau H^{}(I𝒳,)\mathrm{\Lambda }_s`$. This can be found in , Proposition 2.16. ∎
###### Remark 3.1.4.
In untwisted Gromov-Witten theory one usually use the Novikov ring $`\mathrm{\Lambda }_{nov}`$ as the ground ring. Since we will need to compare untwisted theory with twisted theory, we choose to work with the larger ground ring $`\mathrm{\Lambda }_s`$. This only requires minor notational changes applied to discussions in Section 2.5.
### 3.2. Givental’s formalism for twisted theory
The formalism for twisted theory requires a twisted version of the pairing on $`H^{}(I𝒳,)`$ which we call the $`(𝐜,F)`$-twisted orbifold Poincaré pairing $`(,)_{(𝐜,F)}`$. It is defined by
$$(a,b)_{(𝐜,F)}:=_{𝒳_i}aI_i^{}b𝐜((q^{}F)_i^{inv}),\text{for }aH^{}(𝒳_i,),bH^{}(𝒳_{i^I},).$$
For other choices of $`a,b`$ the pairing $`(a,b)_{(𝐜,F)}`$ is defined to be $`0`$.
Consider another symplectic vector space $`(_{(𝐜,F)},\mathrm{\Omega }_{(𝐜,F)})`$, where $`_{(𝐜,F)}=`$ and the $`\mathrm{\Lambda }_s`$-valued symplectic form $`\mathrm{\Omega }_{(𝐜,F)}`$ is given by
$$\mathrm{\Omega }(f,g)_{(𝐜,F)}=\text{Res}_{z=0}(f(z),g(z))_{(𝐜,F)}dz,$$
for $`f,g_{(𝐜,F)}`$.
###### Lemma 3.2.1.
The symplectic vector spaces $`(,\mathrm{\Omega })`$ and $`(_{(𝐜,F)},\mathrm{\Omega }_{(𝐜,F)})`$ are identified via the map
(3.2.1.1)
$$(_{(𝐜,F)},\mathrm{\Omega }_{(𝐜,F)})(,\mathrm{\Omega })$$
defined by $`aa\sqrt{𝐜((q^{}F)^{inv})},`$ where $`a\sqrt{𝐜((q^{}F)^{inv})}`$ is the ordinary cup product in $`H^{}(I𝒳,)`$.
###### Proof.
For $`a,bH^{}(I𝒳)`$, we have
$$\begin{array}{cc}& (a\sqrt{𝐜((q^{}F)^{inv})},b\sqrt{𝐜((q^{}F)^{inv})})_{orb}=_{I𝒳}a\sqrt{𝐜((q^{}F)^{inv})}I^{}(b\sqrt{𝐜((q^{}F)^{inv})})\hfill \\ & =_{I𝒳}a\sqrt{𝐜((q^{}F)^{inv})}I^{}bI^{}\sqrt{𝐜((q^{}F)^{inv})}=_{I𝒳}a\sqrt{𝐜((q^{}F)^{inv})}I^{}b\sqrt{𝐜(I^{}((q^{}F)^{inv}))}\hfill \\ & =_{I𝒳}a\sqrt{𝐜((q^{}F)^{inv})}I^{}b\sqrt{𝐜((q^{}F)^{inv})}=_{I𝒳}aI^{}b𝐜((q^{}F)^{inv})=(a,b)_{(𝐜,F)}.\hfill \end{array}$$
Here the fact $`I^{}((q^{}F)^{inv})=(q^{}F)^{inv}`$ is used, see Lemma 2.2.1 (2). ∎
We equip $`_{(𝐜,F)}`$ with the same polarization as that of $``$, namely $`_{(𝐜,F)}=(_{(𝐜,F)})_+(_{(𝐜,F)})_{}`$ with $`(_{(𝐜,F)})_\pm =_\pm `$. This polarization also identifies $`_{(𝐜,F)}`$ with $`(_{(𝐜,F)})_+(_{(𝐜,F)})_+^{}`$, where $`(_{(𝐜,F)})_+^{}`$ is the dual $`\mathrm{\Lambda }_s`$-module. (We may thus think $`_{(𝐜,F)}`$ as the cotangent bundle $`T^{}(_{(𝐜,F)})_+`$.)
We define the twisted dilaton shift to be $`𝐪(z)=\sqrt{𝐜((q^{}F)^{inv})}(𝐭(z)\mathrm{𝟏}z)`$, where the ordinary cup product in $`H^{}(I𝒳,)`$ is used. Via the twisted dilaton shift the twisted total descendant potential $`𝒟_{(𝐜,F)}(𝐭)`$ is regarded as an element in the Fock space, the space of $`\mathrm{\Lambda }_s[[\mathrm{},\mathrm{}^1]]`$-valued formal functions on $`_+`$ in the formal neighborhood of $`𝐪=\sqrt{𝐜((q^{}F)^{inv})}\mathrm{𝟏}z`$.
Similar to the untwisted case, the twisted genus-$`0`$ potential $`_{(𝐜,F)}^0`$, which is defined in a formal neighborhood of $`\mathrm{𝟏}z_+`$, defines a (formal germ of) Lagrangian submanifold
$$_{(𝐜,F)}:=\{(𝐩,𝐪)|𝐩=d_𝐪_{(𝐜,F)}^0\}.$$
Here $`_{(𝐜,F)}^0`$ is first regarded as an element in the Fock space of functions on $`(_{(𝐜,F)})_+(_{(𝐜,F)},\mathrm{\Omega }_{(𝐜,F)})`$ via the untwisted dilaton shift. Define a (formal germ of) Lagrangian submanifold $`\stackrel{~}{}_{(𝐜,F)}(_{(𝐜,F)},\mathrm{\Omega }_{(𝐜,F)})`$ by the graph of its differential. Second, the map (3.2.1.1) identifies this Lagrangian submanifold with the submanifold $`_{(𝐜,F)}(,\mathrm{\Omega })`$.
We remark that it is not a priori clear whether the Lagrangian submanifold $`_{(𝐜,F)}`$ satisfies (3.1.1.1) or not. This will be a consequence of our main theorem, see Corollary 4.2.3.
###### Definition 3.2.2.
The twisted $`J`$-function $`J_{(𝐜,F)}(t,z)`$ is defined as follows:
$$\begin{array}{cc}\hfill (J_{(𝐜,F)}(t,z),a)_{(𝐜,F)}& :=(z+t,a)_{(𝐜,F)}+\underset{n0,d\text{Eff}(𝒳)}{}\frac{Q^d}{n!}t,\mathrm{},t,\frac{a}{z\overline{\psi }}_{0,n+1,d,(𝐜,F)}\hfill \\ & =(z+t,a)_{(𝐜,F)}+\underset{n0,d\text{Eff}(𝒳)}{}\underset{k0}{}\frac{Q^d}{n!}t,\mathrm{},t,a\overline{\psi }^k_{0,n+1,d,(𝐜,F)}\frac{1}{z^{k+1}}.\hfill \end{array}$$
Again, the twisted $`J`$-function is a formal power series in coordinates $`t^\alpha `$ of $`t=_\alpha t^\alpha \varphi _\alpha H^{}(I𝒳,)`$ taking values in $`_{(𝐜,F)}`$.
### 3.3. Quantization of Quadratic Hamiltonians
Givental observed that many interesting relations in Gromov-Witten theory be expressed in simple forms by applying the Weyl quantization, which is a standard way to produce (projective) Fock space representations of the Heisenberg Lie algebra, to his symplectic space formalism. In this section, we describe this quantization of quadratic Hamiltonian procedure. This quantization procedure allows us to write the quantum Riemann-Roch formula in a simple form.
Let $`A:`$ be a linear infinitesimally symplectic transformation, i.e. $`\mathrm{\Omega }(Af,g)+\mathrm{\Omega }(f,Ag)=0`$ for all $`f,g`$. When $`f`$ is written in Darboux coordinates, the quadratic Hamiltonian
$$f\frac{1}{2}\mathrm{\Omega }(Af,f),$$
is a series of homogeneous degree two monomials in Darboux coordinates $`p_a^\alpha ,q_b^\alpha `$. Define the quantization of quadratic monomials as
$$\widehat{q_a^\mu q_b^\nu }=\frac{q_a^\mu q_b^\nu }{\mathrm{}},\widehat{q_a^\mu p_b^\nu }=q_a^\mu \frac{}{q_b^\nu },\widehat{p_a^\mu p_b^\nu }=\mathrm{}\frac{}{q_a^\mu }\frac{}{q_b^\nu }.$$
Extending linearly, this defines a quadratic differential operator $`\widehat{A}`$, called the quantization of $`A`$. The differential operators $`\widehat{q_aq_b},\widehat{q_ap_b},\widehat{p_ap_b}`$ act on Fock. Since the quadratic Hamiltonian of $`A`$ may contain infinitely many monomials, the quantization $`\widehat{A}`$ do not act on $`\mathrm{𝐹𝑜𝑐𝑘}`$ in general. The quantization of a symplectic transformation of the form $`\mathrm{exp}(A)`$, with $`A`$ infinitesimally symplectic, is defined to be $`\mathrm{exp}(\widehat{A})=_{k0}\frac{\widehat{A}^k}{k!}`$. In general, $`\mathrm{exp}(\widehat{A})`$ is not well-defined. However the operator that occurs in our quantum Riemann-Roch formula does act on the descendant potential.
For infinitesimal symplectomorphisms $`A`$ and $`B`$, there is the following relation
$$[\widehat{A},\widehat{B}]=\{A,B\}^{}+𝒞(h_A,h_B),$$
where $`\{,\}`$ is the Lie bracket, $`[,]`$ is the supercommutator, and $`h_A`$ (respectively $`h_B`$) is the quadratic Hamiltonian of $`A`$ (respectively $`B`$). A direct calculation shows that the cocycle $`𝒞`$ is given by
$$\begin{array}{cc}& 𝒞(p_a^\mu p_b^\nu ,q_a^\mu q_b^\nu )=𝒞(q_a^\mu q_b^\nu ,p_a^\mu p_b^\nu )=1+\delta ^{\mu \nu }\delta _{ab},\hfill \\ & 𝒞=0\text{ on any other pair of quadratic Darboux monomials.}\hfill \end{array}$$
For simplicity, we write $`𝒞(A,B)`$ for $`𝒞(h_A,h_B)`$.
Some universal equations in orbifold Gromov-Witten theory can be expressed as differential equations satisfied by the total descendant potential $`𝒟_𝒳`$. These differential equations can often be written in very simple forms using the quantization formalism. For example,
###### Lemma 3.3.1.
The string equation can be written as
(3.3.1.1)
$$\widehat{\left(\frac{1}{z}\right)}𝒟_𝒳=0.$$
###### Proof.
This is proved in the same way as that for varieties (which can be found in , Example 1.3.3.2). We explain the details for the readers’ convenience.
Put $`𝐭_i(z)=_{j0}t_{ij}z^jH^{}(I𝒳)[[z]]`$. The string equation in cases $`(g,n,d)(0,3,0),(1,1,0)`$ can be written as
$$𝐭_1(\overline{\psi }),\mathrm{},𝐭_{n1}(\overline{\psi }),\mathrm{𝟏}_{g,n,d}=\underset{i=1}{\overset{n1}{}}𝐭_1(\overline{\psi }),\mathrm{},\left[\frac{𝐭_i(\overline{\psi })}{\overline{\psi }}\right]_+,\mathrm{},𝐭_{n1}(\overline{\psi })_{g,n1,d},$$
where
$$\left[\frac{𝐭_i(\overline{\psi })}{\overline{\psi }}\right]_+=\underset{j1}{}t_{ij}\overline{\psi }^{j1}.$$
Summing over $`g,n,d`$ yields
$$\begin{array}{cc}& \underset{g,n,d}{}\frac{Q^d\mathrm{}^{g1}}{(n1)!}𝐭(\overline{\psi }),\mathrm{},𝐭(\overline{\psi }),\mathrm{𝟏}_{g,n,d}\hfill \\ & =\underset{g,n,d}{}\frac{Q^d\mathrm{}^{g1}}{(n1)!}\left[\frac{𝐭(\overline{\psi })}{\overline{\psi }}\right]_+,𝐭(\overline{\psi }),\mathrm{},𝐭(\overline{\psi })_{g,n,d}+\frac{1}{2\mathrm{}}𝐭(\overline{\psi }),𝐭(\overline{\psi }),\mathrm{𝟏}_{0,3,0}+\mathrm{𝟏}_{0,1,0}.\hfill \end{array}$$
This gives
$$\frac{1}{2\mathrm{}}q_0^\alpha g_{\alpha \beta }q_0^\beta \underset{g,n,d}{}\frac{Q^d\mathrm{}^{g1}}{(n1)!}\left[\frac{𝐪(\overline{\psi })}{\overline{\psi }}\right]_+,𝐭(\overline{\psi }),\mathrm{},𝐭(\overline{\psi })_{g,n,d}=0,$$
where $`g_{\alpha \beta }=(\varphi _\alpha ,\varphi _\beta )_{orb}`$.
A direct calculation shows that this is (3.3.1.1). ∎
In the proof of Theorem 4.2.1 we will encounter quantizations of operators of the form $`A=Bz^m`$ with $`B\text{End}(H^{}(I𝒳))`$. An explicit expression of $`\widehat{A}`$ may be found by a straightforward computation. This is worked out in , Example 1.3.3.1, to which we refer the readers for details. See also Appendix C.
### 3.4. Loop space interpretation
In this Section we sketch an interpretation of Givental’s formalism in terms of loop spaces. The interpretation is topological in nature, so we work with the topological stack underlying the Deligne-Mumford stack $`𝒳`$ (which we still denote by $`𝒳`$). We should point out that while this interpretation sheds some light on the conceptual meaning of this formalism, one can work with the formalism without knowing this loop space interpretation.
Let $`L𝒳=Map(S^1,𝒳)`$ be the stack of loops in $`𝒳`$. Definition and properties of $`L𝒳`$ can be found in e.g . Loop rotation yields an $`S^1`$-action on $`L𝒳`$. The stack $`L𝒳^{S^1}`$ of $`S^1`$-fixed loops is identified with the inertia stack $`I𝒳`$. One may think of $`_+`$ as the $`S^1`$-equivariant cohomology of $`L𝒳`$ expressed in terms of the cohomology of the space $`L𝒳^{S^1}I𝒳`$ and the first Chern class $`z`$ of the universal line bundle $`L`$ over $`BS^1`$.
## 4. Quantum Riemann-Roch
As in Section 2.5.8, consider a characteristic class $`𝐜`$ which is multiplicative and invertible. Since the logarithm of $`𝐜`$ is additive, it is a linear combination of components of the Chern character. Hence we may write
$$𝐜()=\mathrm{exp}\left(\underset{k0}{}s_kch_k()\right).$$
For convenience, set $`s_1=0`$. We regard $`s_k`$ as parameters and consider the twisted descendant potentials $`𝒟_s:=𝒟_{(𝐜,F)}`$ as a family of elements in the Fock space of functions on $`_+`$ depending<sup>12</sup><sup>12</sup>12The $`s_k`$-dependence of $`𝒟_{(𝐜,F)}`$ occurs in two places: the twisting factor $`𝐜(F_{g,n,d})`$ and the twisted dilaton shift. on $`s_k`$. We have $`𝒟_s=𝒟_𝒳`$ when all $`s_k=0`$. In this Section we formulate our main result, Theorem 4.2.1, which expresses $`𝒟_s`$ in terms of $`𝒟_𝒳`$.
### 4.1. Some Infinitesimal Symplectic Operators
In this section we introduce certain operators acting on $``$ which will be used in the subsequent sections.
Recall that the Bernoulli polynomials $`B_m(x)`$ are defined by
$$\frac{te^{tx}}{e^t1}=\underset{m0}{}\frac{B_m(x)t^m}{m!},$$
see for instance , Section 7.20. In particular we have $`B_0(x)=1,B_1(x)=x1/2`$. The Bernoulli numbers $`B_m`$ are given by $`B_m:=B_m(0)`$. The following Lemma is immediate from the definition.
###### Lemma 4.1.1.
$`B_m(1x)=(1)^mB_m(x)`$.
###### Definition 4.1.2.
For each integer $`m0`$, define an element $`A_mH^{}(I𝒳)=_iH^{}(𝒳_i)`$ as follows: The component of $`A_m`$ on $`H^{}(𝒳_i)`$ is
$$A_m|_{𝒳_i}:=\underset{0lr_i1}{}ch(F_i^{(l)})B_m(l/r_i).$$
Let $`(A_m)_k`$ denote the degree $`2k`$ part of $`A_m`$,
$$(A_m)_k|_{𝒳_i}:=\underset{0lr_i1}{}ch_k(F_i^{(l)})B_m(l/r_i).$$
Ordinary multiplication by $`A_m`$ defines an operator on $`H^{}(I𝒳)`$. By abuse of notation, we denote this operator by $`A_m`$. The quantization of the operator $`A_mz^{m1}`$ will appear in Theorem 4.2.1. The main goal of this section is to prove that $`A_mz^{m1}`$ is infinitesimally symplectic, which is not a priori clear. It follows from the following result.
###### Lemma 4.1.3.
For $`m1`$, the operator $`A_{2m+1}`$ is anti-self-adjoint with respect to the usual or twisted orbifold Poincaré pairing. The operator $`A_{2m}`$ is self-adjoint with respect to the usual or twisted orbifold Poincaré pairing.
###### Proof.
We prove the statements for the usual pairing. The proofs for the twisted pairing are identical.
For $`aH^{}(𝒳_i),bH^{}(𝒳_{i^I})`$ and $`0<l<r_i`$, we have, by Lemma 2.2.1 (1), the following:
$$(ch(F_i^{(l)})a,b)_{orb}=_{𝒳_i}ch(F_i^{(l)})aI^{}b=_{𝒳_i}aI^{}ch(F_{i^I}^{(r_il)})I^{}b=(a,ch(F_{i^I}^{(r_il)})b)_{orb}.$$
Multiplying by $`B_{2m+1}(l/r_i)`$ yields
$$(B_{2m+1}(l/r_i)ch(F_i^{(l)})a,b)_{orb}=(a,B_{2m+1}(l/r_i)ch(F_{i^I}^{(r_il)})b)_{orb}\text{for }0<l<r_i.$$
By Lemma 4.1.1, $`B_{2m+1}(\frac{l}{r_i})=B_{2m+1}(1\frac{l}{r_i})=B_{2m+1}(\frac{r_il}{r_i})`$. Hence for $`0<l<r_i`$ we have
(4.1.3.1)
$$(B_{2m+1}(l/r_i)ch(F_i^{(l)})a,b)_{orb}=(a,B_{2m+1}(\frac{r_il}{r_i})ch(F_{i^I}^{(r_il)})b)_{orb}.$$
By Lemma 2.2.1 (2), $`(ch(F_i^{(0)})a,b)_{orb}=(a,ch(F_{i^I}^{(0)})b)_{orb}`$. Since $`B_{2m+1}(0)=0`$ for $`m1`$, we have
(4.1.3.2)
$$(B_{2m+1}(0)ch(F_i^{(0)})a,b)_{orb}=(a,B_{2m+1}(0)ch(F_{i^I}^{(0)})b)_{orb}.$$
Adding (4.1.3.1) for $`l=1,\mathrm{},r_i1`$ and (4.1.3.2) yields
$$(A_{2m+1}|_{𝒳_i}a,b)_{orb}=(a,A_{2m+1}|_{𝒳_{i^I}}b)_{orb},$$
which proves the statement about $`A_{2m+1}`$.
To prove the statement for $`A_{2m}`$, we start with
$$(B_{2m}(l/r_i)ch(F_i^{(l)})a,b)_{orb}=(a,B_{2m}(l/r_i)ch(F_{i^I}^{(r_il)})b)_{orb}\text{for }0<l<r_i,$$
and
(4.1.3.3)
$$(B_{2m}(0)ch(F_i^{(0)})a,b)_{orb}=(a,B_{2m}(0)ch(F_{i^I}^{(0)})b)_{orb}.$$
By Lemma 4.1.1, $`B_{2m}(l/r_i)=B_{2m}(1l/r_i)=B_{2m}(\frac{r_il}{r_i})`$. This implies that, for $`0<l<r_i`$,
(4.1.3.4)
$$(B_{2m}(l/r_i)ch(F_i^{(l)})a,b)_{orb}=(a,B_{2m}(\frac{r_il}{r_i})ch(F_{i^I}^{(r_il)})b)_{orb}.$$
Adding (4.1.3.4) for $`l=1,\mathrm{},r_i1`$ and (4.1.3.3) yields
$$(A_{2m}|_{𝒳_i}a,b)_{orb}=(a,A_{2m}|_{𝒳_{i^I}}b)_{orb},$$
which proves the statement about $`A_{2m}`$. ∎
###### Remark 4.1.4.
1. Since $`B_0(x)=1`$, we have
$$A_0|_{𝒳_i}=\underset{0lr_i1}{}ch(F_i^{(l)})=ch(q^{}F)|_{𝒳_i}.$$
Thus multiplication by $`A_0`$ defines a self-adjoint operator with respect to both pairings.
2. The operator $`A_1`$ is not anti-self-adjoint. However note that
$$A_1|_{𝒳_i}=\underset{0lr_i1}{}B_1(l/r_i)ch(F_i^{(l)})=B_1(0)ch(F_i^{(0)})+\underset{l=1}{\overset{r_i1}{}}B_1(l/r_i)ch(F_i^{(l)}).$$
We can use the arguments in the proof of Lemma 4.1.3 to show that
$$(\underset{l=1}{\overset{r_i1}{}}B_1(l/r_i)ch(F_i^{(l)})a,b)_{orb}=(a,\underset{l=1}{\overset{r_i1}{}}B_1(1l/r_i)ch(F_{i^I}^{(r_il)})b)_{orb}.$$
Using $`B_1(0)=1/2`$ we rewrite this as
$$((A_1|_{𝒳_i}+\frac{1}{2}ch(F_i^{(0)}))a,b)_{orb}=(a,(A_1|_{𝒳_{i^I}}+\frac{1}{2}ch(F_{i^I}^{(0)}))b)_{orb}.$$
In our notation, $`ch(F_i^{(0)})=ch((q^{}F)^{inv})|_{𝒳_i}`$. Thus multiplication by $`A_1+\frac{1}{2}ch((q^{}F)^{inv})`$ defines an anti-self-adjoint operator.
3. Lemma 4.1.3 also holds if we replace $`A_{2m+1}`$ and $`A_{2m}`$ by $`(A_{2m+1})_k`$ and $`(A_{2m})_k`$ respectively.
###### Corollary 4.1.5.
Multiplications by the following classes define infinitesimally symplectic transformations on $`(,\mathrm{\Omega })`$ and $`(_{(𝐜,F)},\mathrm{\Omega }_{(𝐜,F)})`$:
$$\begin{array}{cc}& A_{2m}z^{2m1},A_{2m+1}z^{2m},m1;A_0/z,A_1+\frac{1}{2}ch((q^{}F)^{inv});\hfill \\ & (A_{2m})_kz^{2m1},(A_{2m+1})_kz^{2m},m1;(A_0)_k/z,(A_1)_k+\frac{1}{2}ch_k((q^{}F)^{inv}).\hfill \end{array}$$
### 4.2. Orbifold Quantum Riemann-Roch formula
Recall that in the definition of twisted orbifold Gromov-Witten invariants in Section 2.5.8, we assume Assumption 2.5.9 (i.e $`𝒳`$ is assumed to be a quotient of a smooth quasi-projective scheme by a linear algebraic group). This assumption is needed also for the application of Grothendieck-Riemann-Roch. To apply Grothendieck-Riemann-Roch formula for Deligne-Mumford stacks to the universal family of orbifold stable maps, we need the universal family to have certain properties. The required properties are proved in for those $`𝒳`$ that satisfy Assumption 2.5.9. Many interesting stacks, for instance the toric Deligne-Mumford stacks , satisfy Assumption 2.5.9. Through the collective efforts of many works, including , , , it is now known that if $`𝒳`$ is a smooth, separated, generically tame Deligne-Mumford stack over $``$ with quasi-projective coarse moduli space, then $`𝒳`$ satisfies Assumption 2.5.9. See , Section 4 for a detailed account.
Now we state the orbifold quantum Riemann-Roch theorem. Its proof is deferred to Section 7.
###### Theorem 4.2.1 (Orbifold Quantum Riemann-Roch).
Let $`𝒳`$ be as in Assumption 2.5.9. Then we have
$$\begin{array}{cc}& \mathrm{exp}\left(\frac{s_0}{2}\text{rank}F\overline{\psi }_{1,1,0}+s_0c_1(F)_{1,1,0}\right)𝒟_s\hfill \\ & =\mathrm{exp}\left(\underset{k0}{}s_k\left(\underset{m>0}{}\frac{(A_m)_{k+1m}z^{m1}}{m!}+\frac{ch_k((q^{}F)^{inv})}{2}\right)^{}\right)\mathrm{exp}\left(\underset{k0}{}s_k\left(\frac{(A_0)_{k+1}}{z}\right)^{}\right)𝒟_𝒳.\hfill \end{array}$$
This theorem expresses, in a rather nontrivial way, the twisted descendant potential $`𝒟_s`$ in terms of the untwisted potential $`𝒟_𝒳`$.
###### Remark 4.2.2.
The right-hand side of Theorem 4.2.1 is well-defined. The verification of this is a straightforward modification of , Proposition A.0.2 (and the fact that $`\mathrm{\Lambda }_s`$ is equipped with a topology). We omit the details.
Passing to the quasi-classical limit $`\mathrm{}0`$, we find that applying the operator $`\mathrm{exp}(\widehat{A})`$ to $`𝒟_𝒳`$ corresponds to transforming the Lagrangian submanifold $`_𝒳`$ by the (unquantized) operator $`\mathrm{exp}(A)`$. Hence we have the following
###### Corollary 4.2.3.
The Lagrangian submanifolds $`_s:=_{(𝐜,F)}`$ and $`_𝒳`$ are related by
$$\begin{array}{cc}\hfill _s& =\mathrm{exp}\left(\underset{k0}{}s_k\left(\underset{m+h=k+1;m,h0}{}\frac{(A_m)_hz^{m1}}{m!}+\frac{ch_k((q^{}F)^{inv})}{2}\right)\right)_𝒳\hfill \\ & =\mathrm{exp}\left(\underset{m,h0}{}s_{m+h1}\frac{(A_m)_hz^{m1}}{m!}+\underset{k0}{}s_k\frac{ch_k((q^{}F)^{inv})}{2}\right)_𝒳.\hfill \end{array}$$
In particular, $`_s`$ is the germ of a Lagrangian cone swept out by a finite dimensional family of subspaces (i.e. (3.1.1.1) holds for $`_s`$).
When $`𝒳`$ is a variety, the inertia stack $`I𝒳`$ is just $`𝒳`$ itself and Theorem 4.2.1 reduces to , Theorem 1 of Coates-Givental. An interesting feature of Theorem 4.2.1 is the presence of values of Bernoulli polynomials (see the definition of elements $`A_m`$) in place of Bernoulli numbers which appear in the quantum Riemann-Roch theorem for varieties (, Theorem 1). It would be interesting to find a conceptual way to explain why this is the case.
###### Remark 4.2.4 (Loop space interpretation).
There is a heuristic interpretation of the operator $`\mathrm{\Delta }=\mathrm{exp}(_{k0}s_k(_{m0}\frac{(A_m)_{k+1m}z^{m1}}{m!}+\frac{ch_k((q^{}F)^{inv})}{2}))`$ in terms of loop space $`L𝒳`$ (Section 3.4), which we sketch below. On each component $`𝒳_iI𝒳L𝒳^{S^1}`$, the $`S^1`$-action on $`𝒳_i`$ is trivial. This action is related to the $`S^1`$-action on the coarse space $`X_i`$ via the $`r_i`$-fold cover $`S^1S^1`$. We have $`𝒳_i\times _{S^1}ES^1𝒳_i\times BS^1`$. Let $`pr_1,pr_2`$ be the projections to factors and let $`L^{1/r_i}`$ denote the pullback by $`pr_2`$ of the universal line bundle over $`BS^1`$. Define $``$ to be the $`S^1`$-equivariant vector bundle over $`I𝒳\times BS^1`$ whose restriction to $`𝒳_i\times BS^1`$ is $`_{0lr_i1}pr_1^{}F_i^{(l)}(L^{1/r_i})^l`$.
Consider the infinite product
$$\sqrt{𝐜(F^{(0)})}\underset{m=1}{\overset{\mathrm{}}{}}𝐜\left(L^m\right).$$
We interpret this as follows: Let $`s(x):=_{k0}s_k\frac{x^k}{k!}`$. Note that if $`x=c_1(𝔏)`$ is the first Chern class of a line bundle $`𝔏`$, then $`s(x)=_ks_k\mathrm{ch}_k(𝔏)=\mathrm{log}𝐜(𝔏)`$. We write
$$\underset{m>0}{}s\left(x+\frac{l}{r}zmz\right)=\frac{e^{\frac{l}{r}z\frac{}{x}}z\frac{}{x}}{e^{z\frac{}{x}}1}\left(z\frac{}{x}\right)^1s(x)=\underset{m0}{}\frac{B_m(\frac{l}{r})}{m!}\left(z\frac{}{x}\right)^{m1}s(x).$$
Using this (and splitting principle) we expand
$$\mathrm{log}\left(\underset{m>0}{}𝐜(F_i^{(l)}L^{\frac{l}{r_i}m})\right)=\underset{k0}{}s_k\underset{m0}{}\frac{B_m(l/r_i)}{m!}\mathrm{ch}_{k+1m}(F_i^{(l)})z^{m1}.$$
Thus the infinite product above gives rise the operator $`\mathrm{\Delta }`$ after some simplification.
### 4.3. Relations to Hurwitz-Hodge integrals
Let $`GSL_n()`$ be a finite subgroup. Consider a $`G`$-action on $`^n`$ without trivial factors so that $`0^n`$ is an isolated $`G`$-fixed point. Hurwitz-Hodge integrals (c.f. ) arise in the study of orbifold Gromov-Witten theory of $`[^n/G]`$. More precisely, the components of $`\overline{}_{g,n}([^n/G],0)`$ parametrizing maps with images $`[0/G]`$ from orbicurves with stacky marked points may be identified with $`\overline{}_{g,n}(BG)`$, and the restriction of the virtual fundamental class is given by the Euler class $`e(R^1f_{}ev_{n+1}^{}𝒱)`$, where $`𝒱`$ is the vector bundle over $`BG`$ defined by the $`G`$-representation $`^n`$, $`f:\overline{}_{g,n+1}(BG)^{}\overline{}_{g,n}(BG)`$ is the universal orbicurve, and $`ev_{n+1}:\overline{}_{g,n+1}(BG)^{}BG`$ is the universal orbifold stable map (see Section 2.4). The integrals over $`\overline{}_{g,n}(BG)`$ of cohomology classes involving $`e(R^1f_{}ev_{n+1}^{}𝒱)`$ are called Hurwitz-Hodge integrals.
One may consider an equivariant version of this: Let $`^{}`$ acts on $`^n`$ by scaling. This $`^{}`$-action commutes with the $`G`$-action, hence descends to a $`^{}`$-action on the stack $`[^n/G]`$. A $`^{}`$-equivariant Hurwitz-Hodge integral
$$_{[\overline{}_{g,n}(BG)]}(\mathrm{}.)e_{^{}}(R^1f_{}ev_{n+1}^{}𝒱)$$
coincides with
$$e_{^{}}(R^0f_{}ev_{n+1}^{}𝒱)_{[\overline{}_{g,n}(BG)]}(\mathrm{})e_{^{}}^1(𝒱_{g,n,0}),$$
where $`(\mathrm{})`$ denotes cohomology and/or descendant insertions. From this it is easy to conclude that Hurwitz-Hodge integrals can be determined by twisted orbifold Gromov-Witten invariants of $`BG`$. Theorem 4.2.1 implies that descendant twisted orbifold Gromov-Witten invariants of $`BG`$ are determined by the usual descendant orbifold Gromov-Witten invariants of $`BG`$. In , explicit formulas expressing descendant orbifold Gromov-Witten invariants of $`BG`$ in terms of descendant integrals on moduli stacks of stable curves have been proven. It is interesting to combine these results to obtain formulas for Hurwitz-Hodge integrals. In genus zero, under additional assumptions, a procedure of explicitly computing $`(e_{^{}}^1,𝒱)`$-twisted orbifold Gromov-Witten invariants using information about usual orbifold Gromov-Witten invariants of $`BG`$ has been established, see . A method of computing Hurwitz-Hodge integrals directly using the Grothedieck-Riemann-Roch calculation in this paper has been developed by J. Zhou , and is used by him to prove the crepant resolution conjecture in higher genus for type $`A`$ surface singularities .
## 5. Quantum Lefschetz
### 5.1. Twisting by Euler class
Consider the group $`^{}`$ which acts trivially on $`𝒳`$ and on the vector bundle $`F`$ by scaling the fibers. Let $`\lambda `$ be the equivariant parameter and $`e()`$ the $`^{}`$-equivariant Euler class. In this section we consider the special case of twisting by $`F`$ and $`𝐜=e`$. From $`\lambda +x=\mathrm{exp}(_{k0}s_k\frac{x^k}{k!})`$, we find<sup>13</sup><sup>13</sup>13Here we work over the ground ring $`\mathrm{\Lambda }_s`$ with the values of $`s_k`$ specified by (5.1.0.1).
(5.1.0.1)
$$s_k=\{\begin{array}{cc}\hfill \mathrm{ln}\lambda ,& \hfill k=0\\ \hfill \frac{(1)^{k1}(k1)!}{\lambda ^k},& \hfill k>0.\end{array}$$
Let $`\rho _i^{l,j}`$ be the Chern roots of $`F_i^{(l)}`$, $`j=1,\mathrm{},rankF_i^{(l)}`$. The following is the case $`𝐜=e`$ of Corollary 4.2.3.
###### Corollary 5.1.1.
The Lagrangian cone $`_e:=_{(e,F)}`$ of the twisted theory is obtained from $`_𝒳`$ by (ordinary) multiplication by the product over Chern roots $`\rho _i^{l,j}`$ of
$$\gamma _{\rho _i^{l,j}}(z)=\{\begin{array}{cc}\hfill \mathrm{exp}(\frac{(\rho _i^{l,j}+\lambda )\mathrm{ln}(\rho _i^{l,j}+\lambda )(\rho _i^{l,j}+\lambda )}{z}+\mathrm{ln}\lambda (\frac{l}{r_i}\frac{1}{2})+(\frac{l}{r_i}\frac{1}{2})\mathrm{ln}(1+\frac{\rho _i^{l,j}}{\lambda })& \\ \hfill +_{m2}\frac{(1)^mB_m(l/r_i)}{m(m1)}(\frac{z}{\lambda +\rho _i^{l,j}})^{m1}),& \\ \hfill \text{if }l0;& \\ \hfill \mathrm{exp}(\frac{(\rho _i^{0,j}+\lambda )\mathrm{ln}(\rho _i^{0,j}+\lambda )(\rho _i^{0,j}+\lambda )}{z}+_{m2}\frac{(1)^mB_m}{m(m1)}(\frac{z}{\lambda +\rho _i^{0,j}})^{m1}),& \\ \hfill \text{if }l=0.& \end{array}$$
###### Proof.
We substitute the definition of $`s_k`$ into the statement of Corollary 4.2.3 and express components $`ch_h(F_i^{(l)})`$ of the Chern characters using the Chern roots $`\rho _i^{l,j}`$. Then by using the identity
$$\underset{h0}{}s_{m+h1}\frac{\rho ^h}{h!}=\frac{d^{m1}}{d\rho ^{m1}}\mathrm{ln}(\lambda +\rho )=\frac{(1)^m(m2)!}{(\lambda +\rho )^{m1}},\text{for }m1$$
we check directly that the $`z^{m1}`$ terms, with $`m1`$, coincide with what are given in Corollary 4.2.3. For the $`z^1`$ term, a direct calculation gives
$$\frac{1}{z}\underset{k0}{}s_kch_{k+1}(F_i^{(l)})=\frac{1}{z}\underset{j}{}\left[((\rho _i^{l,j}+\lambda )\mathrm{ln}(\rho _i^{l,j}+\lambda )(\rho _i^{l,j}+\lambda ))(\lambda \mathrm{ln}\lambda \lambda )\right].$$
Since the operator $`\frac{1}{z}`$ preserves the cone $`_𝒳`$, we may discard the term $`\frac{\lambda \mathrm{ln}\lambda \lambda }{z}`$. The result follows. ∎
Our next goal is to extract from Corollary 5.1.1 more explicit information about genus zero invariants. For the rest of this section and Section 5.2, we make the following assumption.
###### Assumption 5.1.2.
1. The generic stabilizer of the stack $`𝒳`$ is trivial.
2. The bundle $`F`$ is a direct sum $`_jF_j`$ of line bundles and each $`F_j`$ is a line bundle pulled back via the natural map $`\pi :𝒳X`$ to the coarse moduli space $`X`$.
###### Remark 5.1.3.
1. In the situation of Assumption 5.1.2, the intersection index $`c_1(F_j),\pi ^{}\beta :=c_1(F_j)\pi ^{}\beta `$ is an integer for all effective curve classes $`\beta `$ of $`X`$. Let $`L=\pi ^{}M`$ be a line bundle on $`𝒳`$ that is pulled back from a line bundle $`M`$ on the coarse moduli space $`X`$. Then for any such $`\beta `$, we have $`c_1(L)\pi ^{}\beta =c_1(M)\beta `$.
2. For each $`i`$ and $`j`$, the line bundle $`q^{}(F_j)|_{𝒳_i}`$ has $`\mu _{r_i}`$-eigenvalue $`1`$. In other words, $`q^{}(F_j)|_{𝒳_i}=q^{}(F_j)|_{𝒳_i}^{(0)}`$.
We are interested in a more precise relationship between the $`J`$-function $`J_𝒳`$ and the twisted $`J`$-function $`J_{(e,F)}`$. We generalize the approach of .
###### Definition 5.1.4.
Put $`\rho _{ji}:=c_1(q^{}(F_j)|_{𝒳_i})H^2(𝒳_i)`$ and $`\rho _j:=c_1(F_j)H^2(𝒳)`$. Define $`I_F(t,z):=(I_F(t,z)_i)`$ where
$$I_F(t,z)_i:=\underset{d\text{Eff}(𝒳)}{}J_{i,d}(t,z)Q^d\underset{j}{}\frac{_{k=\mathrm{}}^{\rho _j,d}(\lambda +\rho _{ji}+kz)}{_{k=\mathrm{}}^0(\lambda +\rho _{ji}+kz)}.$$
Following , we call this the hypergeometric modification of $`J_𝒳`$.
###### Remark 5.1.5.
Assumption 5.1.2 is used to ensure that the intersection indices $`c_1(F_j),d`$ are integers, which is needed in order for the hypergeometric modification to be well-defined. (Note that for $`d\text{Eff}(𝒳)`$ there exists $`\beta \text{Eff}(X)`$ such that $`d=\pi ^{}\beta `$. If $`c_1(F_j),d=c_1(F_j),\pi ^{}\beta `$ are not integers, the product $`_{k=\mathrm{}}^{\rho _j,d}(\lambda +\rho _{ji}+kz)/_{k=\mathrm{}}^0(\lambda +\rho _{ji}+kz)`$ doesn’t make sense.)
###### Theorem 5.1.6.
The family
$$tI_F(t,z),tH^{}(I𝒳)$$
of vectors in $`(_{(e,F)},\mathrm{\Omega }_{(e,F)})`$ lies on the Lagrangian submanifold $`\stackrel{~}{}_{(e,F)}`$.
###### Remark 5.1.7.
Theorem 5.1.6 uses Assumption 5.1.2 in a essential way. A more general result of this kind is given in .
###### Proof.
This is a generalization of , Theorem 2 (see also , Theorem 1.7.3). In view of Assumption 5.1.2 and Lemma 2.3.8, we may rewrite the operators $`\gamma _{\rho _j^{l,i}}(z)`$ in terms of the Chen-Ruan orbifold cup product (note that our assumption forces $`l=0`$). More precisely,
$$\gamma _{\rho _j^{0,i}}(z)=\mathrm{exp}(\frac{(\rho _j+\lambda )\mathrm{ln}(\rho _j+\lambda )(\rho _j+\lambda )}{z}+\underset{m2}{}\frac{(1)^mB_m}{m(m1)}\left(\frac{z}{\lambda +\rho _j}\right)^{m1})_{orb}|_{H^{}(𝒳_i)}.$$
It is then straightforward to check that the argument of and applies verbatim (of course with Corollary 5.1.1 replacing its manifold version). Details are left to the readers.
###### Corollary 5.1.8.
The tangent space $`L_t`$ to $`\stackrel{~}{}_{(e,F)}`$ at the point $`I_F(t,z)`$ is equal to the tangent space of $`_{(e,F)}`$ at a unique point $`J_{(e,F)}(\tau (t),z)`$, where $`\tau (t)H^{}(I𝒳,)\mathrm{\Lambda }_s`$.
###### Proof.
Note that $`I_F(t,z)J_𝒳(t,z)modQ`$. An easy calculation shows that the family
$$tI_F(t,z),tH^{}(I𝒳,)\mathrm{\Lambda }_s$$
is transverse to $`zL_t`$ for every $`t`$. As pointed out in Corollary 4.2.3, (3.1.1.1) holds for $`_{(e,F)}`$. Thus the proof of , Proposition 2.16 may be applied to show that $`L_t`$ is equal to the tangent space of $`_{(e,F)}`$ at a unique point $`J_{(e,F)}(\tau (t),z)`$. ∎
###### Remark 5.1.9.
1. Intuitively this Corollary may be interpreted as saying that the intersection of $`zL_t`$ with $`\{z+z_{}\}\stackrel{~}{}_{(e,F)}`$ is equal to
$$J_{(e,F)}(\tau (t),z)z+\tau (t)+_{},$$
where $`\tau (t)H^{}(I𝒳,)\mathrm{\Lambda }_s`$ is defined by this intersection.
2. This Corollary should be viewed as a procedure of computing $`J_{(e,F)}`$ from $`I_F`$. This procedure is related to Birkhoff factorization in the theory of loop groups. More precisely, this procedure applied to the first derivatives of $`I_F`$ is indeed an example of Birkhoff factorization.
3. The map $`t\tau =\tau (t)`$ may be viewed as the “mirror map”. This Corollary gives a geometric description of this map.
### 5.2. Complete intersections
In this Section we apply Corollary 5.1.8 to vector bundles with some positivity property to deduce relationships between orbifold Gromov-Witten invariants of a complete intersection orbifold and orbifold Gromov-Witten invariants of its ambient orbifold.
###### Definition 5.2.1.
A line bundle $`F`$ over $`𝒳`$ is called convex if $`H^1(𝒞,f^{}F)=0`$ for all $`1`$-pointed genus-$`0`$ orbifold stable maps $`f:(𝒞,\mathrm{\Sigma })𝒳`$.
###### Example 5.2.2.
Let $`L:=\pi ^{}M`$ be a line bundle on $`𝒳`$ that is the pullback of a line bundle $`M`$ on the coarse moduli space $`X`$. For an orbifold stable map $`f:𝒞𝒳`$ with induced map $`\overline{f}:CX`$ between coarse moduli spaces, we have
$$H^1(𝒞,f^{}L)=H^1(𝒞,f^{}\pi ^{}M)=H^1(𝒞,\overline{\pi }^{}\overline{f}^{}M)=H^1(C,\overline{f}^{}M).$$
Here $`\overline{\pi }:𝒞C`$ is the map to the coarse curve. From this we see that the bundle $`L`$ is convex if $`M`$ is convex in the usual sense.
The following Proposition follows from .
###### Proposition 5.2.3.
Let $`F=_jF_j`$ be a direct sum of convex line bundles. Let $`𝒴`$ be the zero locus of a regular section of $`F`$, and $`j_{0,n,d}:\overline{}_{0,n}(𝒴,d)\overline{}_{0,n}(𝒳,d)`$ the induced map. Then $`j_{0,n,d}[\overline{}_{0,n}(𝒴,d)]^w=𝔢(F_{0,n,d})[\overline{}_{0,n}(𝒳,d)]^w`$, where $`𝔢()`$ denotes the non-equivariant Euler class.
In the situation of Proposition 5.2.3 let $`j:𝒴𝒳`$ be the inclusion. Let $`I_{𝒳,𝒴}(t,z)`$ and $`J_{𝒳,𝒴}(\tau ,z)`$ be the nonequivariant limits $`\lambda 0`$ of $`I_F(t,z)`$ and $`J_{(e,F)}(\tau ,z)`$ respectively. Let $`F_{0,n+1,d}^{}`$ be the kernel of the evaluation map $`F_{0,n+1,d}ev_{n+1}^{}q^{}F`$ at the $`(n+1)`$-st marked point. Note that the image of the evaluation map is contained in $`ev_{n+1}^{}((q^{}F)^{inv})`$. The non-equivariant limit $`J_{𝒳,𝒴}`$ can be written as
$$J_{𝒳,𝒴}(t,z)=z+t+\underset{n0,d\text{Eff}(𝒳)}{}\frac{Q^d}{n!}ev_{n+1}(ev_1^{}t\mathrm{}ev_n^{}t\frac{𝔢(F_{0,n+1,d}^{})}{z\overline{\psi }_{n+1}}).$$
Together with Proposition 5.2.3 this implies that
(5.2.3.1)
$$𝔢((q^{}F)^{inv})J_{𝒳,𝒴}(u,z)=j_{}J_𝒴(j^{}u,z)$$
where on the right-hand side the Novikov rings should be changed according to $`\text{Eff}(𝒴)\text{Eff}(𝒳)`$.
By taking the nonequivariant limit, we obtain
###### Corollary 5.2.4.
Let $`𝒳,𝒴`$ and $`F`$ be as in Proposition 5.2.3. Then $`I_{𝒳,𝒴}(t,z)`$ and $`J_{𝒳,𝒴}(\tau ,z)`$ determine the same cone. Moreover, $`J_{𝒳,𝒴}(\tau ,z)`$ is determined from $`I_{𝒳,𝒴}(t,z)`$ by the procedure described in Corollary 5.1.8, followed by the mirror map $`t\tau `$.
This is a generalization of “Quantum Lefschetz Hyperplane Principle” (see , , , , , ) to Deligne-Mumford stacks.
We now restrict to the small parameter space $`H^2(𝒳)`$. We continue to assume that $`F=_jF_j`$ is a direct sum of convex line bundles.
###### Proposition 5.2.5.
Let $`\{\gamma _k\}`$ be a basis for $`H^2(𝒳)`$. If $`c_1(F)c_1(T_𝒳)`$, then for $`tH^2(𝒳)`$ we have an expansion
$$I_{𝒳,𝒴}(t,z)=zF(t)+\underset{k}{}G^k(t)\gamma _k+O(z^1),$$
where $`F(t)`$ and $`G^k(t)`$ are certain scalar-valued functions with $`F(t)`$ invertible.
###### Proof.
We have
$$I_F(t,z)_i=z+t+\underset{d>0}{}J_{i,d}(t,z)Q^d\underset{j}{}\underset{k=1}{\overset{\rho _j,d}{}}(\lambda +\rho _{ji}+kz)+O(z^1).$$
Recall that
$$J_{i,d}(t,z)=\underset{n0}{}\frac{Q^d}{n!}\underset{k0,\alpha }{}t,\mathrm{},t,\varphi _\alpha \overline{\psi }^k_{0,n+1,d}\frac{\varphi ^\alpha }{z^{k+1}},$$
where $`\{\varphi ^\alpha \}`$ is an additive basis of $`H^{}(𝒳_i)`$. We need to identify the highest power of $`z`$ in $`J_{i,d}(t,z)`$. For this one should take $`tH^2(𝒳)`$ and $`orbdeg(\varphi _\alpha )`$ to be as large as possible. In view of Lemma 2.3.4, the largest possible orbifold degree is $`2dim_{}𝒳`$. Therefore, by (2.5.2.1), the largest power of $`z`$ in $`J_{i,d}(t,z)`$ is $`1c_1(T𝒳),\pi ^{}d`$.
The highest power of $`z`$ occurring in
$$J_{i,d}(t,z)Q^d\underset{j}{}\underset{k=1}{\overset{\rho _j,d}{}}(\lambda +\rho _{ji}+kz)$$
is equal to
$$1+c_1(F),\pi ^{}dc_1(T_𝒳),\pi ^{}d.$$
By our assumption, this is at most $`1`$. If this is equal to $`1`$, then the class $`\varphi ^\alpha `$ has orbifold degree $`0`$. In order to have $`z^0`$ term, we must have $`orbdeg(\varphi _\alpha )2dim_{}𝒳2`$, which implies that $`orbdeg(\varphi ^\alpha )2`$. Also, we see that $`F(t)1(modQ)`$. The Proposition follows. ∎
Since $`J_{𝒳,𝒴}(\tau ,z)`$ is characterized by the asymptotic $`J_{𝒳,𝒴}(\tau ,z)=z+\tau +O(z^1),`$ by comparing the asymptotics of $`I_{𝒳,𝒴}`$ and $`J_{𝒳,𝒴}`$, we obtain
###### Corollary 5.2.6.
If $`c_1(F)c_1(T_𝒳)`$, then the restriction of $`J_{𝒳,𝒴}(\tau ,z)`$ to small parameter space $`H^2(𝒳)`$ is given by
$$J_{𝒳,𝒴}(\tau ,z)=\frac{I_{𝒳,𝒴}(t,z)}{F(t)},\text{where}\tau =\underset{k}{}\frac{G^k(t)}{F(t)}\gamma _k.$$
This may be regarded as a mirror formula for complete intersection orbifolds. Once the $`J`$-function of $`𝒳`$ is known, part of the $`J`$-function of $`𝒴`$ that involves classes pulled back from $`𝒳`$ can be computed by Corollary 5.2.6 and (5.2.3.1).
## 6. Quantum Serre duality
The so-called “Quantum Serre duality” (, ) is formulated as a relation between $`(𝐜,F)`$-twisted invariants and invariants twisted by the “dual data” $`(𝐜^{},F^{})`$ defined below. In this Section we prove such a relation for Deligne-Mumford stacks.
### 6.1. General case
We again consider the general case of twisting by a vector bundle $`F`$ over $`𝒳`$ and multiplicative invertible characteristic class $`𝐜()=\mathrm{exp}(_ks_kch_k())`$. Here we do not require Assumption 5.1.2. Consider the dual case of twisting by the dual bundle $`F^{}`$ and the class
$$𝐜^{}():=\mathrm{exp}\left(\underset{k0}{}(1)^{k+1}s_kch_k()\right).$$
Note that $`𝐜^{}(F^{})=\frac{1}{𝐜(F)}`$. An application of Theorem 4.2.1 yields the following relation between the potentials $`𝒟_{(𝐜,F)}`$ and $`𝒟_{(𝐜^{},F^{})}`$.
###### Theorem 6.1.1 (Quantum Serre duality for orbifolds).
Let $`𝐭^{}(z)=𝐜((q^{}F)^{inv})𝐭(z)+(1𝐜((q^{}F)^{inv}))z`$. Then we have
$$𝒟_{(𝐜^{},F^{})}(𝐭^{})=\mathrm{exp}\left(s_0\text{rank}F\overline{\psi }_{1,1,0}\right)𝒟_{(𝐜,F)}(𝐭).$$
###### Proof.
One may prove this result by comparing the formulas for $`𝒟_{(𝐜,F)}`$ and $`𝒟_{(𝐜^{},F^{})}`$ given by Theorem 4.2.1. We proceed differently by comparing the differential equation (7.1.1.4) for $`(𝐜,F)`$ and $`(𝐜^{},F^{})`$. The equation satisfied by $`𝒟_{(𝐜,F)}`$ is
(6.1.1.1)
$$\frac{𝒟_{(𝐜,F)}}{s_k}=\left[\left(\underset{m+h=k+1;m,h0}{}\frac{(A_m)_hz^{m1}}{m!}+\frac{ch_k(q^{}F^{inv})}{2}\right)^{}+C_k\right]𝒟_{(𝐜,F)}.$$
We write the equation satisfied by $`𝒟_{(𝐜^{},F^{})}`$ as
(6.1.1.2)
$$(1)^{k+1}\frac{𝒟_{(𝐜^{},F^{})}}{s_k}=\left[\left(\underset{m+h=k+1;m,h0}{}\frac{(A_m^{})_hz^{m1}}{m!}+\frac{ch_k(q^{}F^{inv})}{2}\right)^{}+C_k^{}\right]𝒟_{(𝐜^{},F^{})}.$$
For a fixed $`i`$, the first term on the right-hand side of (6.1.1.2) is
$$\underset{\stackrel{m+h=k+1}{m,h0}}{}\frac{1}{m!}\underset{0lr_i1}{}ch_h(F_i^{(l)})B_m(l/r_i)z^{m1}.$$
We now analyze this for each fixed $`m,h`$. Using $`F_i^{(l)}=F_i^{(r_il)}`$ for $`0<l<r_i`$ and $`F_i^{(0)}=F_i^{(0)}`$, we can write this as
$$\begin{array}{cc}& \frac{1}{m!}\underset{1lr_i1}{}ch_h(F_i^{(r_il)})B_m(l/r_i)z^{m1}+\frac{1}{m!}ch_h(F_i^{(0)})B_mz^{m1}\hfill \\ \hfill =& (1)^{m+h}\frac{1}{m!}\underset{1lr_i1}{}ch_h(F_i^{(r_il)})B_m(\frac{r_il}{r_i})z^{m1}+(1)^h\frac{1}{m!}ch_h(F_i^{(0)})B_mz^{m1}\hfill \\ \hfill =& (1)^{m+h}\frac{1}{m!}\underset{1lr_i1}{}ch_h(F_i^{(l)})B_m(\frac{l}{r_i})z^{m1}+(1)^h\frac{1}{m!}ch_h(F_i^{(0)})B_mz^{m1}.\hfill \end{array}$$
For $`m1`$, since $`(1)^mB_m=B_m`$, this sum is
$$(1)^{k+1}\frac{1}{m!}\underset{0lr_i1}{}ch_h(F_i^{(l)})B_m(\frac{l}{r_i})z^{m1},$$
where we use $`m+h=k+1`$. For $`m=1`$ and $`h=k`$, we have
$$(1)^kch_k(F_i^{(0)})B_1=\frac{1}{2}(1)^kch_k(F_i^{(0)}),$$
which cancels with the term $`ch_k(F_i^{(0)})/2`$.
Therefore we conclude that (6.1.1.2) is
(6.1.1.3)
$$\frac{𝒟_{(𝐜^{},F^{})}}{s_k}=\left[\left(\underset{m+h=k+1;m,h0}{}\frac{(A_m)_hz^{m1}}{m!}+\frac{ch_k(q^{}F^{inv})}{2}\right)^{}+C_k^{}\right]𝒟_{(𝐜^{},F^{})}.$$
By Lemma 7.2.1, $`C_k^{}=0`$ for $`k1`$, and
$$C_0^{}=\frac{1}{2}\text{rank}F^{}\overline{\psi }_{1,1,0}c_1(F^{})_{1,1,0}=\frac{1}{2}\text{rank}F\overline{\psi }_{1,1,0}+c_1(F)_{1,1,0}.$$
The result follows by comparing (6.1.1.3) with (6.1.1.1), ∎
### 6.2. Euler class
We consider the case of twisting by a $`^{}`$-equivariant Euler class $`e()`$, where $`^{}`$ acts on $`F`$ by scaling the fibers. Let the dual bundle $`F^{}`$ be equipped with the dual $`^{}`$-action and let $`e^1()`$ be the inverse $`^{}`$-equivariant Euler class. If $`\rho _j`$ are the Chern roots of $`F`$, then $`e^1(F^{})=_j(\lambda \rho _j)^1`$. The main result of this section, Proposition 6.2.1, is a relation between $`(e,F)`$-twisted invariants and $`(e^1,F)`$-twisted invariants. Note that this is not a special case of Theorem 6.1.1 since $`e^1e^{}`$.
Let $`F`$ be a vector bundle on $`𝒳`$. For a component $`𝒳_i`$ of $`I𝒳`$ we define the age of the bundle $`F`$ on $`𝒳_i`$ to be
$$age(F_i):=\underset{1lr_i1}{}\frac{l}{r_i}rankF_i^{(l)}.$$
The bundle $`F_i^{mov}`$ is defined to be $`_{1lr_i1}F_i^{(l)}`$. Let $`M:H^{}(I𝒳)H^{}(I𝒳)`$ be defined as follows: on $`H^{}(𝒳_i)`$, $`M`$ is the multiplication by the number $`(1)^{\frac{1}{2}rankF_i^{mov}age(F_i)}`$. Put
$$𝐭^{}(z)=z+(1)^{\frac{1}{2}rank(q^{}F)^{inv}}Me((q^{}F)^{inv})(𝐭(z)1z),$$
and define a change $`:Q^dQ^d(1)^{ch_1(F),d}`$ in the Novikov ring.
###### Proposition 6.2.1.
We have
$$\begin{array}{cc}& \mathrm{exp}\left(\frac{\pi \sqrt{1}}{2}\text{rank}F\overline{\psi }_{1,1,0}+\pi \sqrt{1}c_1(F)_{1,1,0}\right)𝒟_{(e^1,F^{})}(𝐭^{},Q)\hfill \\ & =\mathrm{exp}\left(\mathrm{ln}\lambda \text{rank}F\overline{\psi }_{1,1,0}\right)𝒟_{(e,F)}(𝐭,Q).\hfill \end{array}$$
###### Proof.
If we write $`e^1()=\mathrm{exp}(_{k0}s_k^{}ch_k())`$ and $`e()=\mathrm{exp}(_{k0}s_kch_k())`$, then we find that $`s_k^{}=(1)^{k+1}s_k`$ for $`k>0`$ and $`s_0^{}=s_0\pi \sqrt{1}`$. The proof of Theorem 6.1.1 shows that $`𝒟_s^{}`$ satisfies the differential equation
(6.2.1.1)
$$\frac{𝒟_s^{}}{s_k}=\left[\left(\underset{m+h=k+1;m,h0}{}\frac{(A_m)_hz^{m1}}{m!}+\frac{ch_k(q^{}F^{inv})}{2}\right)^{}+C_k^{}\right]𝒟_s^{}.$$
Also,
$$𝒟_s^{}|_{s_k=0}=𝒟_s^{}|_{s_0^{}=\pi \sqrt{1},s_k^{}=0\text{ for }k>0}.$$
By Theorem 4.2.1, we have
$$\begin{array}{cc}& \mathrm{exp}\left(\frac{\pi \sqrt{1}}{2}\text{rank}F^{}\overline{\psi }_{1,1,0}\pi \sqrt{1}c_1(F^{})_{1,1,0}\right)𝒟_s^{}|_{\stackrel{s_0^{}=\pi \sqrt{1}}{s_k^{}=0\text{ for }k>0}}\hfill \\ \hfill =& \mathrm{exp}(\pi \sqrt{1}[(A_1)_0+\frac{1}{2}ch_0(q^{}F^{inv})]^{})\mathrm{exp}(\pi \sqrt{1}((A_0)_1/z)^{})𝒟_0.\hfill \end{array}$$
For a fixed $`i`$, we have
$$\begin{array}{cc}& ((A_1)_0+\frac{1}{2}ch_0(q^{}F^{inv}))|_{𝒳_i}=\underset{0<l<r_i}{}ch_0(F_i^{(l)})B_1(l/r_i),\hfill \\ & (A_0)_1/z|_{𝒳_i}=ch_1(q^{}F|_{𝒳_i})/z.\hfill \end{array}$$
The operator
$$\mathrm{exp}(\pi \sqrt{1}((A_0)_1/z)^{})$$
can be computed directly using Appendix C and is seen to yield the change $``$ via the divisor flow. The operator $`\mathrm{exp}(\pi \sqrt{1}[(A_1)_0+\frac{1}{2}ch_0(q^{}F^{inv})]^{})`$ may be computed using Appendix C, one sees that it yields the operator $`M`$.
Solving the equation (6.2.1.1) and using the expression of $`𝒟_s^{}|_{s_k=0}`$ yields the desired formula. ∎
## 7. Proof of Theorem 4.2.1
In this section we prove Theorem 4.2.1. The proof is rather lengthy and somewhat unpleasant, however the idea (which we borrowed from ) of the proof is quite simple.
### 7.1. Overview
For the convenience of what follows, we introduce a new notation.
###### Definition 7.1.1.
Let $`a_jH^{p_j}(𝒳_{i_j},),j=1,\mathrm{},n`$ be cohomology classes, $`AH^{}(\overline{}_{g,n}(𝒳,d;i_1,\mathrm{},i_n))`$, and $`k_1,\mathrm{},k_n`$ nonnegative integers. Define
$$a_1\overline{\psi }^{k_1},\mathrm{},a_n\overline{\psi }^{k_n};A_{g,n,d}:=_{[\overline{}_{g,n}(𝒳,d;i_1,\mathrm{},i_n)]^w}(ev_1^{}a_1)\overline{\psi }_1^{k_1}\mathrm{}(ev_n^{}a_n)\overline{\psi }_n^{k_n}A.$$
Let us explain the structure of the proof. As explained in Section 4, the twisted descendant potentials $`𝒟_s`$ are viewed as a family of asymptotic elements depending on variables $`s=(s_0,s_1,\mathrm{})`$. We know that
$$𝒟_s|_{s_0=s_1=\mathrm{}=0}=𝒟_𝒳.$$
To prove Theorem 4.2.1, we find a system of differential equations in $`s_k`$ satisfied by $`𝒟_s`$, and solve the initial value problem with the initial condition given by above. Such a system of differential equations is found by doing the naive thing: compute $`𝒟_s/s_k`$. A direct computation yields
$$𝒟_s^1\frac{𝒟_s}{s_k}$$
(7.1.1.1)
$$=\underset{g,n,d}{}\frac{Q^d\mathrm{}^{g1}}{n!}𝐭(z),\mathrm{},𝐭(z);\frac{}{s_k}𝐜(F_{g,n,d})_{g,n,d}+\underset{g,n,d}{}\frac{Q^d\mathrm{}^{g1}}{(n1)!}\frac{}{s_k}𝐭(z),\mathrm{},𝐭(z);𝐜(F_{g,n,d})_{g,n,d}.$$
The second term in (7.1.1.1), called the derivative term, is equal to
(7.1.1.2)
$$\frac{1}{2}\underset{g,n,d}{}\frac{Q^d\mathrm{}^{g1}}{(n1)!}ch_k((q^{}F)^{inv})(𝐭(z)1z),\mathrm{},𝐭(z);𝐜(F_{g,n,d})_{g,n,d}.$$
This can be seen from
$$\begin{array}{cc}& \frac{}{s_k}𝐭(z)=\frac{}{s_k}(𝐜((q^{}F)^{inv})^{1/2}𝐪(z)+1z)\hfill \\ & =\frac{1}{2}𝐜((q^{}F)^{inv})^{1/2}ch_k((q^{}F)^{inv})𝐪(z)=\frac{1}{2}ch_k((q^{}F)^{inv})(𝐭(z)1z).\hfill \end{array}$$
Since
$$\frac{}{s_k}𝐜(F_{g,n,d})=𝐜(F_{g,n,d})ch_k(F_{g,n,d}),$$
the first term in (7.1.1.1) is equal to
(7.1.1.3)
$$\underset{g,n,d}{}\frac{Q^d\mathrm{}^{g1}}{n!}𝐭(z),\mathrm{},𝐭(z);𝐜(F_{g,n,d})ch_k(F_{g,n,d})_{g,n,d}.$$
The Chern character $`ch_k(F_{g,n,d})`$ appearing in (7.1.1.3) will be computed by applying Grothendieck-Riemann-Roch formula. The result is then combined with (7.1.1.2) to obtained the following differential equation, written using Givental’s formalism:
(7.1.1.4)
$$\frac{𝒟_s}{s_k}=\left[\left(\underset{m+h=k+1;m,h0}{}\frac{(A_m)_hz^{m1}}{m!}+\frac{ch_k((q^{}F)^{inv})}{2}\right)^{}+C_k\right]𝒟_s.$$
Here we define
$$C_k:=_{[\overline{}_{1,1}(𝒳,0)^{}]^w}\underset{a+b=k+1;a,b0}{}ev^{}ch_a(F)(Td^{}(L_1))_b𝐜(F_{1,1,0}).$$
The term $`()_b`$ means the degree $`2b`$ part of a cohomology class, and $`Td^{}`$ is the dual Todd class defined by the property that $`Td^{}(L^{})=Td(L)`$ for any line bundle $`L`$. Recall that the superscript indicates the quantization, discussed in Section 3.3.
The proof of Theorem 4.2.1 will be completed in the next few sections. In the next Section we derive Theorem 4.2.1 from (7.1.1.4). The computation of $`ch_k(F_{g,n,d})`$ by applying Grothendieck-Riemann-Roch formula will be presented in Section 7.3. In Section 7.4 we derive the equation (7.1.1.4) from these computations.
### 7.2. From (7.1.1.4) to (4.2.1)
We first derive Theorem 4.2.1 from (7.1.1.4).
###### Lemma 7.2.1.
$`C_k=0`$ for $`k1`$. $`C_0=\frac{1}{2}\text{rank}F\overline{\psi }_{1,1,0}c_1(F)_{1,1,0}`$.
###### Proof.
The virtual complex dimension of $`\overline{}_{1,1}(𝒳,0)^{}`$ is $`1`$ (note that the marked point is non-stacky). The integrand involved in $`C_k`$ is of degree at least $`2(k+1)`$. So $`C_k=0,k1`$ for dimension reason. The degree $`2`$ part of the integrand of $`C_0`$ is $`(ev^{}ch_0(F)(Td^{}(L_1))_1+ev^{}ch_1(F))𝐜(F_{1,1,0})_0`$, where $`𝐜(F_{1,1,0})_0=\mathrm{exp}(s_0ch_0(F_{1,1,0}))`$ denotes the degree-$`0`$ part of $`𝐜(F_{1,1,0})`$. By Riemann-Roch, we find that the virtual rank of $`F_{1,1,0}`$ is $`0`$, thus $`ch_0(F_{1,1,0})=0`$ and $`𝐜(F_{1,1,0})_0=1`$. We conclude by observing that $`(Td^{}(L_1))_1=\frac{1}{2}\overline{\psi }`$. ∎
###### Remark 7.2.2.
Our proof of Lemma 7.2.1 uses a dimension argument and is valid in non-equivariant Gromov-Witten theory. The exact evaluation of $`C_k`$ in equivariant Gromov-Witten theory requires an explicit description of the moduli stack $`\overline{}_{1,1}(𝒳,0)^{}`$ and its virtual class. Such a description is not known for Deligne-Mumford stacks $`𝒳`$, thus an exact evaluation of $`C_k`$ in equivariant Gromov-Witten theory remains unknown. If the torus acts with isolated fixed points, virtual localization formula yields a calculation of $`C_k`$. We will not discuss it here.
For simplicity, write
$$\alpha _k:=\left(\underset{m>0}{}\frac{(A_m)_{k+1m}z^{m1}}{m!}+\frac{ch_k(q^{}F^{inv})}{2}\right)^{},\beta _k:=\left(\frac{(A_0)_{k+1}}{z}\right)^{}.$$
As explained in , Example 1.3.4.1, the cocycle $`𝒞(_js_j\alpha _j,\beta _k)`$ is equal to
(7.2.2.1)
$$\begin{array}{cc}& 𝒞(\underset{j0}{}\frac{s_j(A_2)_{j1}z}{2},\frac{(A_0)_{k+1}}{z})\hfill \\ \hfill =& \frac{1}{2}\text{str}\left((A_0)_{k+1}\underset{j0}{}\frac{s_j(A_2)_{j1}}{2}\right)\hfill \\ \hfill =& \frac{1}{2}_{II𝒳}e(T_{II𝒳})(Iq)^{}(A_0)_{k+1}\underset{j>0}{}s_j\frac{(Iq)^{}(A_2)_{j1}}{2}\text{by Appendix }\text{D}\hfill \\ \hfill =& 0\text{since the degrees of integrands exceed the dimension of }II𝒳.\hfill \end{array}$$
Solving (7.1.1.4), we find
$$𝒟_s=\mathrm{exp}(\underset{k}{}s_k\alpha _k)\mathrm{exp}(\underset{k}{}s_k\beta _k)\mathrm{exp}(s_0C_0)𝒟_0,$$
which gives Theorem 4.2.1.
###### Remark 7.2.3.
Our derivation of Theorem 4.2.1 from (7.1.1.4) uses the exact values of $`C_k`$ and the cocycles, and is valid for non-equivariant Gromov-Witten theory. In this paper we only consider non-equivariant Gromov-Witten theory. Note however that (7.1.1.4) is valid in full generality.
### 7.3. GRR Calculation
In this Section we compute $`ch_k(F_{g,n,d})[\overline{}_{g,n}(𝒳,d)]^{vir}`$ by applying Grothendieck-Riemann-Roch formula. For technical reasons we proceed as follows. The construction in using Hilbert functors for Deligne-Mumford stacks provides a family of orbicurves
$$𝒰$$
over a smooth base stack $``$ and an embedding
$$\overline{}_{g,n}(𝒳,d)$$
satisfying the following
###### Property 7.3.1.
1. the family $`𝒰`$ pulls back to the universal family over $`\overline{}_{g,n}(𝒳,d)`$,
2. the vector bundle $`E=ev_{n+1}^{}F`$ extends to a vector bundle over $`𝒰`$,
3. the Kodaira-Spencer map $`T_mExt^1(𝒪_{𝒰_m},𝒪_{𝒰_m})`$ is surjective for all $`m`$.
Details can be found in , Proposition 3.1.1. We check that Grothendieck-Riemann-Roch formula (Corollary A.0.7) can be applied to $`𝒰`$. First note that Property 7.3.1 and the smoothness of $``$ imply that $`𝒰`$ is a smooth Deligne-Mumford stack. By the construction in , $`𝒰`$ factors as
$$𝒰\overline{𝔸}\times ,$$
where $`\overline{𝔸}`$ is smooth, $`𝒰\overline{𝔸}\times `$ is a regular embedding, and $`\overline{𝔸}\times `$ is the projection. Therefore $`𝒰`$ is a local complete intersection (lci) morphism<sup>14</sup><sup>14</sup>14The notion of a lci morphism for Deligne-Mumford stacks is the same as that for schemes (, Appendix B.7.6). . Moreover, since the relative tangent bundle of $`\overline{𝔸}\times `$ is just the tangent bundle of $`\overline{𝔸}`$ pulled back to $`\overline{𝔸}\times `$, it follows that the lci virtual tangent bundle of $`𝒰`$ coincides with its relative tangent bundle.
We can compute $`ch(f_{}ev_{n+1}^{}F)[\overline{}_{g,n}(𝒳,d)]^{vir}`$ by applying Corollary A.0.7 to the morphism $`𝒰`$ then capping with $`[\overline{}_{g,n}(𝒳,d)]^{vir}`$. Therefore for the rest of this section, we assume Property 7.3.1. To avoid introducing cumbersome new notations, we will express our computations as if they were done for the morphism $`\overline{}_{g,n+1}(𝒳,d)^{}\overline{}_{g,n}(𝒳,d)`$. Namely we assume that the moduli stack $`\overline{}_{g,n}(𝒳,d)`$ is smooth and its universal family has everywhere-surjective Kodaira-Spencer map.
The Grothendieck-Riemann-Roch calculation we need is done individually for each component $`\overline{}_{g,n+1}(𝒳,d;i_1,\mathrm{},i_n,0)`$. We begin with an analysis of the components of the inertia stacks $`I\overline{}_{g,n+1}(𝒳,d;i_1,\mathrm{},i_n,0)`$ required for this calculation. There are three types of components of the inertia stack $`I\overline{}_{g,n+1}(𝒳,d;i_1,\mathrm{},i_n,0)`$ that are mapped to $`\overline{}_{g,n}(𝒳,d;i_1,\mathrm{},i_n)`$:
1. the main stratum $`\overline{}_{g,n+1}(𝒳,d;i_1,\mathrm{},i_n,0)`$,
2. the divisors of marked points $`𝒟_{j,(i_1,\mathrm{},i_n)}`$, and
3. the locus of nodes $`𝒵_{r,(i_1,\mathrm{},i_n)}`$.
In the rest of this Section we work out contributions from each of them to GRR formula of $`ch(f_{}ev_{n+1}^{}F)[\overline{}_{g,n}(𝒳,d)]^{vir}`$.
#### 7.3.2. Main stratum
The computation on the main stratum does not depend on $`(i_1,\mathrm{},i_n)`$. To simplify notation, we describe it for $`f:\overline{}_{g,n+1}(𝒳,d)^{}\overline{}_{g,n}(𝒳,d)`$.
The restrictions of $`\stackrel{~}{ch}(E)`$ and $`\stackrel{~}{Td}(T_f)`$ to $`\overline{}_{g,n+1}(𝒳,d)^{}`$ are $`ch(E)`$ and $`Td(T_f)`$ respectively. To compute $`Td(T_f)=Td^{}(\mathrm{\Omega }_f)`$, we use the following Lemma.
###### Lemma 7.3.3.
There are exact sequences of sheaves
$$0\mathrm{\Omega }_f\omega _fi_{}𝒪_𝒵0,$$
$$0\omega _fL_{n+1}_js_j𝒪_{𝒟_j}0,$$
where
* $`L_{n+1}`$ is the tautological line bundle on $`\overline{}_{g,n+1}(𝒳,d)^{}`$ corresponding to the $`(n+1)`$-st marked point.
* $`s_j:𝒟_j\overline{}_{g,n+1}(𝒳,d)^{}`$ are inclusions of the divisors of marked points.
* $`i:𝒵\overline{}_{g,n+1}(𝒳,d)^{}`$ is the inclusion of the locus of the nodes.
###### Proof.
We prove the first exact sequence. The second sequence can be proved by a similar argument. Away from $`𝒵`$, two sheaves $`\mathrm{\Omega }_f`$ and $`\omega _f`$ are the same. Consider a family $`S𝒞\stackrel{𝑓}{}𝒳`$ of orbifold stable maps with $`S=\mathrm{Spec}R`$ such that the fiber of $`𝒞/S`$ over a point of $`S`$ is a nodal orbicurve. Étale-locally near a node<sup>15</sup><sup>15</sup>15We use the condition on Kodaira-Spencer map to give this description., we may write $`𝒞`$ as the quotient $`[U/\mu _r]`$ where $`U`$ is the nodal curve $`\mathrm{Spec}(R[z,w]/(zwt))`$ and $`\mu _r`$ acts on $`U`$ via
$$(z,w)(\zeta _rz,\zeta _r^1w).$$
On this neighborhood, the dualizing sheaf $`\omega _f`$ corresponds to the $`\mu _r`$-equivariant sheaf $`\omega _U`$ with invariant generator $`\frac{dzdw}{d(zw)}`$. The sheaf $`\mathrm{\Omega }_f`$ of Kähler differentials corresponds to the $`\mu _r`$-equivariant sheaf $`\mathrm{\Omega }_U`$ with generators $`dz,dw`$ and a relation $`wdz+zdw=0`$. There is an equivariant inclusion $`\mathrm{\Omega }_f\omega _f`$ defined by
$$dzz\frac{dzdw}{d(zw)},dww\frac{dzdw}{d(zw)}.$$
The cokernel corresponds to the $`\mu _r`$-equivariant sheaf generated by $`\frac{dzdw}{d(zw)}`$ with coefficients in $`𝒪_S`$. This sheaf is identified with $`i_{}𝒪_𝒵`$, proving the first exact sequence. ∎
Therefore we have
$$Td^{}(\mathrm{\Omega }_f)=Td^{}(L_{n+1})Td^{}(i_{}𝒪_𝒵))_jTd^{}(s_j𝒪_{𝒟_j}).$$
Note that the $`𝒟_j`$’s and $`𝒵`$ are disjoint, and the restrictions of $`L_{n+1}`$ to them are trivial. So we have
$$\begin{array}{cc}& (Td^{}(s_{j_1}𝒪_{𝒟_{j_1}})1)(Td^{}(s_{j_2}𝒪_{𝒟_{j_2}})1)=0\text{for }1j_1<j_2n,\hfill \\ & (Td^{}(s_j𝒪_{𝒟_j})1)(Td^{}(L_{n+1})1)=0\text{for }1jn,\hfill \\ & (Td^{}(s_j𝒪_{𝒟_j})1)(Td^{}(i_{}𝒪_𝒵)1)=0\text{for }1jn,\hfill \\ & (Td^{}(i_{}𝒪_𝒵)1)(Td^{}(L_{n+1})1)=0.\hfill \end{array}$$
Equivalently,
$$\begin{array}{cc}& Td^{}(s_{j_1}𝒪_{𝒟_{j_1}}s_{j_2}𝒪_{𝒟_{j_2}})1=(Td^{}(s_{j_1}𝒪_{𝒟_{j_1}})1)+(Td^{}(s_{j_2}𝒪_{𝒟_{j_2}})1)\text{for }1j_1<j_2n,\hfill \\ & Td^{}(s_j𝒪_{𝒟_j}+L_{n+1})1=(Td^{}(s_j𝒪_{𝒟_j})1)+(Td^{}(L_{n+1})1)\text{for }1jn,\hfill \\ & Td^{}(s_j𝒪_{𝒟_j}i_{}𝒪_𝒵)1=(Td^{}(s_j𝒪_{𝒟_j})1)+(Td^{}(i_{}𝒪_𝒵)1)\text{for }1jn,\hfill \\ & Td^{}(i_{}𝒪_𝒵+L_{n+1})1=(Td^{}(i_{}𝒪_𝒵)1)+(Td^{}(L_{n+1})1).\hfill \end{array}$$
Using these equations repeatedly, we find
$$\begin{array}{cc}\hfill Td^{}(\mathrm{\Omega }_f)1& =Td^{}\left(L_{n+1}+\underset{j}{}(s_j𝒪_{𝒟_j})i_{}𝒪_𝒵\right)1\hfill \\ & =(Td^{}(L_{n+1})1)+\underset{j}{}(Td^{}(s_j𝒪_{𝒟_j})^11)+(Td^{}(i_{}𝒪_𝒵)^11).\hfill \end{array}$$
Hence the contribution from the main stratum is
$$\begin{array}{cc}\hfill \{f_{}(ch(E)Td^{}(L_{n+1}))& +\underset{j}{}f_{}(ch(E)(Td^{}(s_j𝒪_{𝒟_j})^11))\hfill \\ & +f_{}(ch(E)(Td^{}(i_{}𝒪_𝒵)^11))\}[\overline{}_{g,n}(𝒳,d)]^{vir}.\hfill \end{array}$$
The term $`Td^{}(s_j𝒪_{𝒟_j})^11`$ is computed as follows: Consider the exact sequence
(7.3.3.1)
$$0𝒪(𝒟_j)𝒪s_j𝒪_{𝒟_j}0.$$
Note that $`s_j^{}(𝒟_j)=c_1(N_j^{})`$ with $`N_j^{}`$ the conormal bundle of $`𝒟_j\overline{}_{g,n+1}(𝒳,d)^{}`$. It follows that
$$\begin{array}{cc}\hfill Td^{}(s_j𝒪_{𝒟_j})^11& =Td^{}(𝒪(𝒟_j))1=\underset{r1}{}\frac{B_r}{r!}(𝒟_j)^r\hfill \\ & =s_j\underset{r1}{}\frac{B_r}{r!}(c_1(N_j^{}))^{r1}=s_j\left[\frac{Td^{}(N_j^{})}{c_1(N_j^{})}\right]_+.\hfill \end{array}$$
Here and henceforth the symbol $`[]_+`$ denotes power series truncation, which removes terms containing negative powers of cohomology classes.
The term $`Td^{}(i_{}𝒪_𝒵)^11`$ is computed as follows: Let $`\varphi :\stackrel{~}{𝒵}𝒵`$ be the double cover of $`𝒵`$ consisting of nodes and choices of a branch at each node. $`\stackrel{~}{𝒵}`$ is a disjoint union of open-and-closed substacks of the form $`\overline{}_{g1,n+\{+,\}}(𝒳,d)\times _{I𝒳\times I𝒳}I𝒳`$ or of the form $`\overline{}_+\times _{I𝒳}\overline{}_{}`$, where $`\overline{}_\pm =\overline{}_{g_\pm ,n_\pm +1}(𝒳,d_\pm )`$ such that $`g_++g_{}=g,n_++n_{}=n,d_++d_{}=d`$ is an ordered splitting of $`g,n,d`$. This follows from the fact that $`𝒵`$ is the universal gerbe of nodes over $`f(𝒵)`$ (c.f. , Proposition 5.2.1).
Let $`L_+`$ be the line bundle on $`\overline{}_+`$ whose fiber at an orbifold stable map is the cotangent space<sup>16</sup><sup>16</sup>16This is not the cotangent space on the coarse curve. at the marked point of gluing. The line bundle $`L_{}`$ on $`\overline{}_{}`$ is similarly defined. On $`\overline{}_{g1,n+\{+,\}}(𝒳,d)`$ the cotangent line bundles at marked points $`+`$ and $``$ are also denoted by $`L_+`$ and $`L_{}`$.
By Lemma 5.1, there is a polynomial $`P`$ such that
$$Td^{}(i_{}𝒪_𝒵)^11=i_{}P(c_1(N),c_2(N)),$$
where $`N`$ is the normal bundle of $`𝒵\overline{}_{g,n+1}(𝒳,d)^{}`$. Thus we have
$$Td^{}(i_{}𝒪_𝒵)^11=\frac{1}{2}i_{}\varphi _{}P(c_1(\varphi ^{}N),c_2(\varphi ^{}N)).$$
Denote $`\iota =i\varphi `$. Using $`\varphi ^{}N=L_+^{}L_{}^{}`$ and the expression of $`P`$ in , page 303, we find
(7.3.3.2)
$$\begin{array}{cc}\hfill Td^{}(i_{}𝒪_𝒵)^11& =\frac{1}{2}\iota _{}\left(\underset{s2}{}\frac{B_s}{s!}\underset{a+b=s2}{}(1)^a\psi _+^a\psi _{}^b\right)\hfill \\ & =\frac{1}{2}\iota _{}\left(\frac{1}{\psi _++\psi _{}}\left(\frac{1}{e^{\psi _+}1}\frac{1}{\psi _+}+\frac{1}{2}+\frac{1}{e^\psi _{}1}\frac{1}{\psi _{}}+\frac{1}{2}\right)\right)\hfill \\ & =\frac{1}{2}\iota _{}\left[\frac{1}{\psi _++\psi _{}}\left(\frac{Td^{}(L_+)}{\psi _+}+\frac{Td^{}(L_{})}{\psi _{}}\right)\right]_+.\hfill \end{array}$$
Here $`\psi _\pm =c_1(L_\pm )`$.
Therefore the contribution from the main stratum is
$$\begin{array}{cc}& f_{}(ch(E)Td^{}(L_{n+1}))[\overline{}_{g,n}(𝒳,d)]^{vir}\hfill \\ & \underset{j}{}f_{}s_j\left(ch(s_j^{}E)\left[\frac{Td^{}(N_j^{})}{c_1(N_j^{})}\right]_+\right)[\overline{}_{g,n}(𝒳,d)]^{vir}\hfill \\ & +\frac{1}{2}(f\iota )_{}\left(ch(\iota ^{}E)\left[\frac{1}{\psi _++\psi _{}}\left(\frac{Td^{}(L_+)}{\psi _+}+\frac{Td^{}(L_{})}{\psi _{}}\right)\right]_+\right)[\overline{}_{g,n}(𝒳,d)]^{vir}.\hfill \end{array}$$
The contribution to $`ch(f_{(i_1,\mathrm{},i_n)}E)[\overline{}_{g,n}(𝒳,d;i_1,\mathrm{},i_n)]^{vir}`$ from the main stratum can be found by restricting the above to $`\overline{}_{g,n}(𝒳,d;i_1,\mathrm{},i_n)`$. It is the sum of the following three terms:
(7.3.3.3)
$$f_{(i_1,\mathrm{},i_n)}(ch(E)Td^{}(L_{n+1}))[\overline{}_{g,n}(𝒳,d;i_1,\mathrm{},i_n)]^{vir},$$
(7.3.3.4)
$$\underset{j}{}f_{(i_1,\mathrm{},i_n)}s_{j,(i_1,\mathrm{},i_n)}\left(ch(s_{j,(i_1,\mathrm{},i_n)}^{}E)\left[\frac{Td^{}(N_j^{})}{c_1(N_j^{})}\right]_+\right)[\overline{}_{g,n}(𝒳,d;i_1,\mathrm{},i_n)]^{vir},$$
(7.3.3.5)
$$\frac{1}{2}(f_{(i_1,\mathrm{},i_n)}\iota _{(i_1,\mathrm{},i_n)})_{}\left(ch(\iota _{(i_1,\mathrm{},i_n)}^{}E)\left[\frac{1}{\psi _++\psi _{}}\left(\frac{Td^{}(L_+)}{\psi _+}+\frac{Td^{}(L_{})}{\psi _{}}\right)\right]_+\right)[\overline{}_{g,n}(𝒳,d;i_1,\mathrm{},i_n)]^{vir}.$$
Here the subscript $`_{(i_1,\mathrm{},i_n)}`$ indicates the restriction to $`\overline{}_{g,n}(𝒳,d;i_1,\mathrm{},i_n)`$. We call (7.3.3.3) the codim-0 term, (7.3.3.4) the codim-1 term, and (7.3.3.5) the codim-2 term.
###### Remark 7.3.4.
Consider the stack $`\overline{}_+\times _{I𝒳}\overline{}_{}`$ parametrizing maps whose domains consist of two parts separated by a distinguished node<sup>17</sup><sup>17</sup>17By definition, a section of the gerbe at the distinguished node is part of the data in this moduli problem.. If the order of the automorphism group of this node is $`r`$, then $`r\psi _+=\overline{\psi }_+`$ is the first Chern class of the line bundle whose fiber is the cotangent line of the coarse curve at the marked point of gluing. Similarly $`r\psi _{}=\overline{\psi }_{}`$. The same statements hold for $`L_\pm `$ on $`\overline{}_{g1,n+\{+,\}}(𝒳,d)`$.
#### 7.3.5. Marked points
We compute the contribution from the divisors formed by marked points. Let $`s_{j,(i_1,\mathrm{},i_n)}:𝒟_{j,(i_1,\mathrm{},i_n)}\overline{}_{g,n+1}(𝒳,d;i_1,\mathrm{},i_n,0)`$ be the divisor of the $`j`$-th marked point. We know that $`𝒟_{j,(i_1,\mathrm{},i_n)}\overline{}_{g,n}(𝒳,d;i_1,\mathrm{},i_n)\times B\mu _{r_{i_j}}`$ and the diagram
$$\begin{array}{ccc}\overline{}_{g,n}(𝒳,d;i_1,\mathrm{},i_n)\times B\mu _{r_{i_j}}& & 𝒳\\ & & \\ \overline{}_{g,n}(𝒳,d;i_1,\mathrm{},i_n).\end{array}$$
defined by the restriction of the universal orbifold stable map is equivalent to the evaluation map
$$ev_j:\overline{}_{g,n}(𝒳,d;i_1,\mathrm{},i_n)𝒳_{i_j}.$$
Also, the generator $`𝔲_{r_{i_j}}\mu _{r_{i_j}}`$ acts on the conormal bundle $`N_j^{}`$ with eigenvalue $`\zeta _{r_{i_j}}^1`$.
The locus $`𝒟_{j,(i_1,\mathrm{},i_n)}`$ contributes components of $`I\overline{}_{g,n+1}(𝒳,d;i_1,\mathrm{},i_n,0)`$ which are mapped to $`\overline{}_{g,n}(𝒳,d;i_1,\mathrm{},i_n)`$. These components are
$$\overline{}_{g,n}(𝒳,d;i_1,\mathrm{},i_n)\times (IB\mu _{r_{i_j}}B\mu _{r_{i_j}})=:\underset{1lr_{i_j}1}{}𝒟_{j,(i_1,\mathrm{},i_n)}(l)$$
where $`𝒟_{j,(i_1,\mathrm{},i_n)}(l)`$ is defined as follows. The inertia stack $`IB\mu _{r_{i_j}}`$ can be described as
$$IB\mu _{r_{i_j}}=\underset{0kr_{i_j}1}{}[\text{Spec}/C(𝔲_{r_{i_j}}^k)].$$
Define
$$𝒟_{j,(i_1,\mathrm{},i_n)}(l):=\overline{}_{g,n}(𝒳,d;i_1,\mathrm{},i_n)\times [\text{Spec}/C(𝔲_{r_{i_j}}^l)]\overline{}_{g,n}(𝒳,d;i_1,\mathrm{},i_n)\times B\mu _{r_{i_j}}.$$
These components arise in the following way. The auotomorphism group of an object of $`𝒟_{j,(i_1,\mathrm{},i_n)}`$ splits as a product $`Aut\times \mu _{r_{i_j}}`$ where the first factor $`Aut`$ is the automorphism group of the correponding object in $`\overline{}_{g,n}(𝒳,d;i_1,\mathrm{},i_n)`$. The components $`𝒟_{j,(i_1,\mathrm{},i_n)}(l),1lr_{i_j}1`$ correspond to taking the identity element of the factor $`Aut`$ and elemenets $`𝔲_{r_{i_j}}^l,1lr_{i_j}1`$ in the second factor $`\mu _{r_{i_j}}`$.
By Lemma 7.3.3 and the exact sequence (7.3.3.1), we see that the pullback of $`T_{f_{(i_1,\mathrm{},i_n)}}`$ to $`𝒟_{j,(i_1,\mathrm{},i_n)}(l)`$ has trivial invariant part, and the moving part is the pullback of $`N_j`$ to $`𝒟_{j,(i_1,\mathrm{},i_n)}(l)`$.
The restriction $`E|_{𝒟_{j,(i_1,\mathrm{},i_n)}}`$ is decomposed into a direct sum $`_{0kr_{i_j}1}E^{(k)}`$ of $`𝔲_{r_{i_j}}`$-eigenbundles, where $`E^{(k)}`$ has $`𝔲_{r_{i_j}}`$-eigenvalue $`\zeta _{r_{i_j}}^k`$ and $`\zeta _{r_{i_j}}=\mathrm{exp}(2\pi \sqrt{1}\frac{1}{r_{i_j}})`$. Let $`P_l:𝒟_{j,(i_1,\mathrm{},i_n)}(l)𝒟_{j,(i_1,\mathrm{},i_n)}`$ be the projection. Then we have
$$ch(\rho (P_l^{}E|_{𝒟_{j,(i_1,\mathrm{},i_n)}}))=\underset{0kr_{i_j}1}{}\zeta _{r_{i_j}}^{kl}ch(P_l^{}(E^{(k)})).$$
So the contribution from $`𝒟_{j,(i_1,\mathrm{},i_n)}(l)`$ is
$$(f_{(i_1,\mathrm{},i_n)}s_{j,(i_1,\mathrm{},i_n)}P_l)_{}\left(\frac{_{0kr_{i_j}1}\zeta _{r_{i_j}}^{kl}ch(P_l^{}(E^{(k)}))}{1\zeta _{r_{i_j}}^lch(P_l^{}N_j^{})}\right)[\overline{}_{g,n}(𝒳,d;i_1,\mathrm{},i_n)]^{vir}.$$
Let $`\gamma _l:\overline{}_{g,n}(𝒳,d;i_1,\mathrm{},i_n)𝒟_{j,(i_1,\mathrm{},i_n)}(l)\overline{}_{g,n}(𝒳,d;i_1,\mathrm{},i_n)\times B\mu _{r_{i_j}}`$ be the map such that the map to the first factor is the identity and the map to the second factor corresponds to the trivial $`\mu _{r_{i_j}}`$-bundle. We have
$$\gamma _l\gamma _l^{}=r_{i_j}id\text{and }f_{(i_1,\mathrm{},i_n)}s_{j,(i_1,\mathrm{},i_n)}P_l\gamma _l=id.$$
Hence we can write the contribution from $`𝒟_{j,(i_1,\mathrm{},i_n)}(l)`$ as
$$\frac{1}{r_{i_j}}\left(\frac{_{0kr_{i_j}1}\zeta _{r_{i_j}}^{kl}ch(\gamma _l^{}P_l^{}(E^{(k)}))}{1\zeta _{r_{i_j}}^lch(\gamma _l^{}P_l^{}N_j^{})}\right)[\overline{}_{g,n}(𝒳,d;i_1,\mathrm{},i_n)]^{vir}.$$
The following Lemma is straightforward.
###### Lemma 7.3.6.
1. For $`E=ev_{n+1}^{}F`$, we have $`\gamma _l^{}P_l^{}(E^{(k)})=ev_j^{}(F_{i_j}^{(k)})`$.
2. $`\gamma _l^{}P_l^{}N_j^{}=L_j`$.
###### Proof.
The second statement follows from the definition. We prove the first statement. Let $`S\overline{}_{g,n}(𝒳,d)`$ be a morphism and $`S𝒞𝒳`$ the corresponding orbifold stable map. Restricting to the divisor of the $`j`$-th marked point yields morphisms
$$S\stackrel{𝑝}{}S\times B\mu _{r_{i_j}}\stackrel{𝜌}{}𝒳.$$
By the description of the inertia stack $`I𝒳`$ in Remark 2.1.2 (i), these morphisms correspond to a morphism $`\stackrel{~}{\rho }:S𝒳_{i_j}`$. Consider the component $`B\mu _{r_{i_j}}[\text{Spec}/C(𝔲_{r_{i_j}}^l)]IB\mu _{r_{i_j}}`$ and let $`\pi _l:S\times [\text{Spec}/C(𝔲_{r_{i_j}}^l)]S\times B\mu _{r_{i_j}}`$ be the projection. Let $`\gamma :SS\times [\text{Spec}/C(𝔲_{r_{i_j}}^l)]`$ be the section of $`p\pi _l`$ such that the map to the first factor is the identity and the map to the second factor corresponds to the trivial $`\mu _{r_{i_j}}`$-bundle. Let $`(\rho ^{}F)^{(k)}`$ be the eigen sub-bundle of $`\rho ^{}F`$ on which $`𝔲_{r_{i_j}}`$ acts with eigenvalue $`\zeta _{r_{i_j}}^k`$. To prove the first statement it suffices to prove
$$\gamma ^{}\pi _l^{}((\rho ^{}F)^{(k)})=\stackrel{~}{\rho }^{}(F_{i_j}^{(k)}).$$
This is obtained immediately from (2.2.1.1) by pulling back via $`\stackrel{~}{\rho }`$.
Therefore the contribution from $`𝒟_{j,(i_1,\mathrm{},i_n)}(l)`$ can be written as
$$\frac{1}{r_{i_j}}\left(\frac{_{0kr_{i_j}1}\zeta _{r_{i_j}}^{kl}ch(ev_j^{}(F_{i_j}^{(k)}))}{1\zeta _{r_{i_j}}^lch(L_j)}\right)[\overline{}_{g,n}(𝒳,d;i_1,\mathrm{},i_n)]^{vir}.$$
The contributions from $`𝒟_{j,(i_1,\mathrm{},i_n)}(1),\mathrm{},𝒟_{j,(i_1,\mathrm{},i_n)}(r_{i_j}1)`$ add up to
$$\begin{array}{cc}& \underset{1lr_{i_j}1}{}\frac{1}{r_{i_j}}\left(\frac{_{0kr_{i_j}1}\zeta _{r_{i_j}}^{kl}ch(ev_j^{}(F_{i_j}^{(k)}))}{1\zeta _{r_{i_j}}^lch(L_j)}\right)[\overline{}_{g,n}(𝒳,d;i_1,\mathrm{},i_n)]^{vir}\hfill \\ & =\frac{1}{r_{i_j}}\underset{0kr_{i_j}1}{}ch(ev_j^{}(F_{i_j}^{(k)}))\underset{1lr_{i_j}1}{}\frac{\zeta _{r_{i_j}}^{kl}}{1\zeta _{r_{i_j}}^le^{c_1(L_j)}}[\overline{}_{g,n}(𝒳,d;i_1,\mathrm{},i_n)]^{vir}.\hfill \end{array}$$
For each $`k`$ with $`0k<r_{i_j}`$ we have
$$\underset{1lr_{i_j}1}{}\frac{\zeta _{r_{i_j}}^{kl}}{1\zeta _{r_{i_j}}^le^{c_1(L_j)}}=\frac{r_{i_j}e^{kc_1(L_j)}}{1e^{r_{i_j}c_1(L_j)}}\frac{1}{1e^{c_1(L_j)}}.$$
Using $`\gamma _0\gamma _0^{}=r_{i_j}id`$ we can rewrite the part of codim-1 term (7.3.3.4) that comes from $`𝒟_{j,(i_1,\mathrm{},i_n)}`$ as
$$\begin{array}{cc}& \frac{1}{r_{i_j}}\underset{0kr_{i_j}1}{}ch(ev_j^{}(F_{i_j}^{(k)}))\underset{n1}{}\frac{B_n}{n!}c_1(L_j)^{n1}[\overline{}_{g,n}(𝒳,d;i_1,\mathrm{},i_n)]^{vir}\hfill \\ & =\frac{1}{r_{i_j}}\underset{0kr_{i_j}1}{}ch(ev_j^{}(F_{i_j}^{(k)}))\frac{1}{c_1(L_j)}\left(\frac{c_1(L_j)}{e^{c_1(L_j)}1}1\right)[\overline{}_{g,n}(𝒳,d;i_1,\mathrm{},i_n)]^{vir}.\hfill \end{array}$$
Here we use $`_{n1}\frac{B_n}{n!}x^{n1}=\frac{1}{x}(\frac{x}{e^x1}1)`$. Combining this part of codim-1 term and contributions from $`𝒟_{j,(i_1,\mathrm{},i_n)}(1)`$,…, $`𝒟_{j,(i_1,\mathrm{},i_n)}(r_{i_j}1)`$, we find that their sum is equal to
$$\frac{1}{r_{i_j}}\underset{0kr_{i_j}1}{}ch(ev_j^{}(F_{i_j}^{(k)}))\left(\frac{r_{i_j}e^{kc_1(L_j)}}{1e^{r_{i_j}c_1(L_j)}}+\frac{1}{c_1(L_j)}\right)[\overline{}_{g,n}(𝒳,d;i_1,\mathrm{},i_n)]^{vir}.$$
Using the definition of Bernoulli polynomials, we see that this is
$$\underset{0kr_{i_j}1}{}ch(ev_j^{}(F_{i_j}^{(k)}))\underset{n1}{}\frac{B_n(k/r_{i_j})}{n!}(r_{i_j}c_1(L_j))^{n1}[\overline{}_{g,n}(𝒳,d;i_1,\mathrm{},i_n)]^{vir}$$
(7.3.6.1)
$$=\underset{n1}{}\frac{_{0kr_{i_j}1}ch(ev_j^{}(F_{i_j}^{(k)}))B_n(k/r_{i_j})}{n!}\overline{\psi }_j^{n1}[\overline{}_{g,n}(𝒳,d;i_1,\mathrm{},i_n)]^{vir}.$$
Here we also use $`\overline{\psi }_j=r_{i_j}\psi _j`$, which follows from the fact that $`L_j^{r_{i_j}}=\pi _n^{}L_j`$.
#### 7.3.7. Nodes
We proceed to compute the contributions from the locus of nodes in a similar fashion. Let
$$\varphi _{r,(i_1,\mathrm{},i_n)}:\stackrel{~}{𝒵}_{r,(i_1,\mathrm{},i_n)}𝒵_{r,(i_1,\mathrm{},i_n)}$$
be the double covering of $`𝒵_{r,(i_1,\mathrm{},i_n)}`$ consisting of nodes and choices of a branch at each node and
$$\iota _{r,(i_1,\mathrm{},i_n)}:\stackrel{~}{𝒵}_{r,(i_1,\mathrm{},i_n)}\overline{}_{g,n+1}(𝒳,d;i_1,\mathrm{},i_n,0)$$
be $`\varphi _{r,(i_1,\mathrm{},i_n)}`$ followed by the inclusion.
The components $`\stackrel{~}{𝒵}_{r,(i_1,\mathrm{},i_n)}(1)`$,…, $`\stackrel{~}{𝒵}_{r,(i_1,\mathrm{},i_n)}(r1)`$ of $`I\stackrel{~}{𝒵}_{r,(i_1,\mathrm{},i_n)}`$ which are mapped to $`\overline{}_{g,n}(𝒳,d;i_1,\mathrm{},i_n)`$ can be defined similarly as $`𝒟_{j,(i_1,\mathrm{},i_n)}(l)`$. Since $`\stackrel{~}{𝒵}`$ can be identified with a disjoint union of the stack
$$\overline{}_{g1,n+\{+,\}}(𝒳,d)\times _{I𝒳\times I𝒳}I𝒳$$
and stacks of the form $`\overline{}_+\times _{I𝒳}\overline{}_{}`$, each $`\stackrel{~}{𝒵}_{r,(i_1,\mathrm{},i_n)}(l)`$ is isomorphic to $`\stackrel{~}{𝒵}_{r,(i_1,\mathrm{},i_n)}`$. Let
$$P_l:\stackrel{~}{𝒵}_{r,(i_1,\mathrm{},i_n)}(l)\stackrel{~}{𝒵}_{r,(i_1,\mathrm{},i_n)}$$
be the projection, and $`\gamma _l`$ an inverse of $`P_l`$. Note that $`\gamma _l\gamma _l^{}=id`$.
By the Koszul complex
$$0𝒪(L_+L_{})𝒪(L_+)𝒪(L_{})𝒪𝒪_{\stackrel{~}{𝒵}}0$$
and Lemma 7.3.3, we see that the invariant part of the pullback of $`T_{f_{(i_1,\mathrm{},i_n)}}`$ to $`\stackrel{~}{𝒵}_{r,(i_1,\mathrm{},i_n)}(l)`$ is the sum of a trivial bundle $`𝒪`$, and $`𝒪`$, and $`(L_+L_{})^{}`$. The moving part is $`(L_+^{}L_{}^{})`$.
The contribution from $`\stackrel{~}{𝒵}_{r,(i_1,\mathrm{},i_n)}(l)`$ is
$$\frac{1}{2}(f_{(i_1,\mathrm{},i_n)}\iota _{r,(i_1,\mathrm{},i_n)})_{}P_l\gamma _l\left(\frac{_k\zeta _r^{lk}ch(\gamma _l^{}P_l^{}(E^{(k)}))Td((L_+L_{})^{})}{1\zeta _r^lch(L_+)\zeta _r^lch(L_{})+ch(^2(L_+L_{}))}\right).$$
Put $`\psi _+=c_1(L_+),\psi _{}=c_1(L_{})`$. Note that $`Td((L_+L_{})^{})=Td^{}((L_+L_{}))=1/Td^{}(L_+L_{})`$. We have
$$\begin{array}{cc}& \frac{Td((L_+L_{})^{})}{1\zeta _r^lch(L_+)\zeta _r^lch(L_{})+ch(^2(L_+L_{}))}\hfill \\ \hfill =& \frac{e^{\psi _++\psi _{}}1}{(\psi _++\psi _{})(1\zeta _r^le^{\psi _+}\zeta _r^le^\psi _{}+e^{\psi _++\psi _{}})}\hfill \\ \hfill =& \frac{e^{\psi _++\psi _{}}1}{(\psi _++\psi _{})(1\zeta _r^le^{\psi _+})(1\zeta _r^le^\psi _{})}=\frac{1}{\psi _++\psi _{}}\left(1+\frac{1}{\zeta _r^le^{\psi _+}1}+\frac{1}{\zeta _r^le^\psi _{}1}\right).\hfill \end{array}$$
Also, for $`0<k<r`$,
$$\underset{l=1}{\overset{r1}{}}\frac{\zeta _r^{kl}}{\zeta _r^le^x1}=\frac{re^{kx}}{e^{rx}1}\frac{1}{e^x1},\underset{l=1}{\overset{r1}{}}\frac{\zeta _r^{kl}}{\zeta _r^le^x1}=\frac{re^{(rk)x}}{e^{rx}1}\frac{1}{e^x1}.$$
And
$$\underset{l=1}{\overset{r1}{}}\frac{1}{\zeta _r^le^x1}=\underset{l=1}{\overset{r1}{}}\frac{1}{\zeta _r^le^x1}=\frac{r}{e^{rx}1}\frac{1}{e^x1},\underset{l=1}{\overset{r1}{}}\zeta _r^{kl}=\{\begin{array}{cc}\hfill 1,& \hfill k0\\ \hfill r1,& \hfill k=0.\end{array}$$
Therefore
(7.3.7.1)
$$\begin{array}{cc}& \underset{l=1}{\overset{r1}{}}\left(\zeta _r^{kl}+\frac{\zeta _r^{kl}}{\zeta _r^le^{\psi _+}1}+\frac{\zeta _r^{kl}}{\zeta _r^le^\psi _{}1}\right)\hfill \\ \hfill =& \{\begin{array}{cc}\hfill \frac{re^{k\psi _+}}{e^{r\psi _+}1}\frac{1}{e^{\psi _+}1}+\frac{re^{(rk)\psi _{}}}{e^{r\psi _{}}1}\frac{1}{e^\psi _{}1}1,& \hfill k0\\ \hfill \frac{r}{e^{r\psi _+}1}\frac{1}{e^{\psi _+}1}+\frac{r}{e^{r\psi _{}}1}\frac{1}{e^\psi _{}1}+r1,& \hfill k=0.\end{array}\hfill \end{array}$$
We have $`\gamma _l^{}P_l^{}E^{(k)}=ev_{node}^{}(F^{(k)})`$, which is similar to Lemma 7.3.6 (see Appendix B for the definition of $`ev_{node}`$). What we need to do now is to combine the part of codim-2 term (7.3.3.5) from $`\stackrel{~}{𝒵}_{r,(i_1,\mathrm{},i_n)}`$ and contributions of $`\stackrel{~}{𝒵}_{r,(i_1,\mathrm{},i_n)}(1),\mathrm{},\stackrel{~}{𝒵}_{r,(i_1,\mathrm{},i_n)}(r1)`$. First note that the term $`ch(\iota _{(i_1,\mathrm{},i_n)}^{}E)`$ in (7.3.3.5) breaks into a sum of terms $`ch(ev_{node}^{}((q^{}F)^{(k)}))`$ for $`0k<r`$. The term in (7.3.3.5) corresponding to $`k`$ is (the pushforward of) $`ch(ev_{node}^{}((q^{}F)^{(k)}))`$ multiplied by $`\frac{1}{\psi _++\psi _{}}`$ and
$$\frac{1}{e^{\psi _+}1}\frac{1}{\psi _+}+\frac{1}{2}+\frac{1}{e^\psi _{}1}\frac{1}{\psi _{}}+\frac{1}{2},$$
see (7.3.3.2). Adding this to (7.3.7.1), we get for $`k0`$
$$\begin{array}{cc}& \frac{re^{k\psi _+}}{e^{r\psi _+}1}\frac{1}{\psi _+}+\frac{re^{(rk)\psi _{}}}{e^{r\psi _{}}1}\frac{1}{\psi _{}}\hfill \\ & =r\left(\frac{e^{\frac{k}{r}r\psi _+}}{e^{r\psi _+}1}\frac{1}{r\psi _+}+\frac{e^{\frac{rk}{r}r\psi _{}}}{e^{r\psi _{}}1}\frac{1}{r\psi _{}}\right)\hfill \\ & =r\underset{n1}{}\left(\frac{B_n(k/r)}{n!}(r\psi _+)^{n1}+\frac{B_n(1k/r)}{n!}(r\psi _{})^{n1}\right),\hfill \end{array}$$
and for $`k=0`$
$$\begin{array}{cc}& \frac{r}{e^{r\psi _+}1}\frac{1}{\psi _+}+\frac{r}{e^{r\psi _{}}1}\frac{1}{\psi _{}}+r\hfill \\ & =r\left(\frac{1}{e^{r\psi _+}1}\frac{1}{r\psi _+}+\frac{1}{2}+\frac{1}{e^{r\psi _{}}1}\frac{1}{r\psi _{}}+\frac{1}{2}\right)\hfill \\ & =r\underset{n2}{}\left(\frac{B_n}{n!}(r\psi _+)^{n1}+\frac{B_n}{n!}(r\psi _{})^{n1}\right)\hfill \\ & =r\underset{n1}{}\left(\frac{B_n}{n!}(r\psi _+)^{n1}+\frac{B_n(1)}{n!}(r\psi _{})^{n1}\right).\hfill \end{array}$$
By these calculation it follows that the combined contribution is the following expression capped with the virtual class:
$$\begin{array}{cc}& \frac{r^2}{2}\left(f_{(i_1,\mathrm{},i_n)}\iota _{r,(i_1,\mathrm{},i_n)}\right)_{}\underset{n1}{}\frac{1}{n!}\frac{1}{r\psi _++r\psi _{}}\underset{0l<r}{}ch\left(ev_{node}^{}\left(\left(q^{}F\right)^{\left(l\right)}\right)\right)\left(B_n\left(\frac{l}{r}\right)\left(r\psi _+\right)^{n1}+B_n\left(1\frac{l}{r}\right)\left(r\psi _{}\right)^{n1}\right)\hfill \\ & =\frac{r^2}{2}(f_{(i_1,\mathrm{},i_n)}\iota _{r,(i_1,\mathrm{},i_n)})_{}\underset{n1}{}\frac{1}{n!}\left[\underset{l}{}ch(ev_{node}^{}((q^{}F)^{(l)}))B_n(l/r)\right]\underset{a+b=n2}{}(1)^b(r\psi _+)^a(r\psi _{})^b\hfill \\ & =\frac{r^2}{2}(f\iota )_{}\underset{n2}{}\frac{1}{n!}\left[\underset{l}{}ch(ev_{node}^{}((q^{}F)^{(l)}))B_n(l/r)\right]\frac{(\overline{\psi }_+)^{n1}+(1)^n(\overline{\psi }_{})^{n1}}{\overline{\psi }_++\overline{\psi }_{}}.\hfill \end{array}$$
Here we use $`r\psi _\pm =\overline{\psi }_\pm `$. Note that we rewrite $`\frac{1}{\psi _++\psi _{}}`$ as $`r\frac{1}{r\psi _++r\psi _{}}`$, which gives a factor of $`r`$.
Combining all together, we find
(7.3.7.2)
$$\begin{array}{cc}& ch(f_{}ev_{n+1}^{}F)[\overline{}_{g,n}(𝒳,d)]^{vir}\hfill \\ & =f_{}(ch(ev^{}F)Td^{}(L_{n+1}))[\overline{}_{g,n}(𝒳,d)]^{vir}\hfill \\ & \underset{i=1}{\overset{n}{}}\underset{m1}{}\frac{ev_i^{}A_m}{m!}(\overline{\psi }_i)^{m1}[\overline{}_{g,n}(𝒳,d)]^{vir}\hfill \\ & +\frac{1}{2}(f\iota )_{}\underset{m2}{}\frac{1}{m!}r_{node}^2(ev_{node}^{}A_m)\left(\frac{(\overline{\psi }_+)^{m1}+(1)^m(\overline{\psi }_{})^{m1}}{\overline{\psi }_++\overline{\psi }_{}}\right)[\overline{}_{g,n}(𝒳,d)]^{vir}.\hfill \end{array}$$
### 7.4. Finding the differential equation
We begin with the following splitting property of the virtual fundamental classes, which will be used in the calculations. Let $`𝔐_{g,n}^{tw}`$ be the (Artin) stack of twisted curves of genus $`g`$ with $`n`$ marked gerbes (not trivialized). First consider the case of separating nodes. Let
$$𝔇^{tw}(g_+;n_+|g_{};n_{}):=\underset{\{1,\mathrm{},n\}=AB,|A|=n_+,|B|=n_{}}{}𝔇^{tw}(g_+;A|g_{};B),$$
where the right-hand side is defined as in , Section 5.1. There is a natural forgetful map $`\overline{}_{g,n}(𝒳,d)𝔐_{g,n}^{tw}`$ and a natural gluing map $`gl:𝔇^{tw}(g_+;n_+|g_{};n_{})𝔐_{g,n}^{tw}`$ as defined in , Proposition 5.1.3. Consider the cartesian diagram formed by these maps:
$$\begin{array}{ccc}𝔇_{g,n}(𝒳)& & \overline{}_{g,n}(𝒳,d)\\ & & & & \\ 𝔇^{tw}(g_+;n_+|g_{};n_{})& \stackrel{gl}{}& 𝔐_{g,n}^{tw}.\end{array}$$
There is a natural map
$$𝔤:\underset{d=d_++d_{}}{}\overline{}_{g_+,n_++1}(𝒳,d_+)\times _{I𝒳}\overline{}_{g_{},n_{}+1}(𝒳,d_{})𝔇_{g,n}(𝒳).$$
This is the universal gerbe over the distinguished node (see , Proposition 5.2.1).
Similarly, for non-separating nodes we write $`gl`$ for the map obtained by gluing the last two marked points. There is a similar cartesian diagram and a similar map $`𝔤`$, which we do not describe explicitly.
###### Proposition 7.4.1.
Let $`\overline{}_{g_+,n_++1}(𝒳,d_+)\times _{I𝒳}\overline{}_{g_{},n_{}+1}(𝒳,d_{})\stackrel{~}{𝒵}_r\stackrel{\iota _r}{}\overline{}_{g,n}(𝒳,d)`$. Consider the diagram of gluing,
$$\begin{array}{ccc}\overline{}_{g_+,n_++1}(𝒳,d_+)\times _{I𝒳}\overline{}_{g_{},n_{}+1}(𝒳,d_{})& & I𝒳\\ & & \delta & & \\ \overline{}_{g_+,n_++1}(𝒳,d_+)\times \overline{}_{g_{},n_{}+1}(𝒳,d_{})& \stackrel{ev_+\times \stackrel{ˇ}{ev}_{}}{}& I𝒳\times I𝒳.\end{array}$$
Here $`\delta :I𝒳I𝒳\times I𝒳`$ is the diagonal map, and $`\stackrel{ˇ}{ev}_{}`$ is the composite
$$\overline{}_{g_{},n_{}+1}(𝒳,d_{})\stackrel{ev_{}}{}I𝒳\stackrel{𝐼}{}I𝒳.$$
Then
$$\underset{d_++d_{}=d}{}\delta ^!([\overline{}_{g_+,n_++1}(𝒳,d_+)]^w\times [\overline{}_{g_{},n_{}+1}(𝒳,d_{})]^w)=r^2𝔤^{}(gl^![\overline{}_{g,n}(𝒳,d)]^w).$$
Similarly, for $`\overline{}_{g1,n+\{+,\}}(𝒳,d)\times _{I𝒳\times I𝒳}I𝒳\stackrel{~}{𝒵}_r\stackrel{\iota _r}{}\overline{}_{g,n}(𝒳,d)`$, we have
$$\delta ^![\overline{}_{g1,n+\{+,\}}(𝒳,d)]^w=r^2𝔤^{}(gl^![\overline{}_{g,n}(𝒳,d)]^w).$$
This Proposition is more general than Proposition 5.3.1 of . The proof of this Proposition is the same as that of , Proposition 5.3.1, with the straightforward adjustment for weighted virtual classes. In particular, the factor $`r^2`$ arises since when a stacky node of order $`r`$ is split into two stacky marked points, each marked point should receive a factor of $`r`$ in order to get the weighted virtual class. (Note that $`r`$ should be interpreted as a locally constant function.)
We now process the term (7.1.1.3). According to the GRR calculation (7.3.7.2), (7.1.1.3) splits into three parts:
Codim-1:
(7.4.1.1)
$$\underset{g,n,d}{}\frac{Q^d\mathrm{}^{g1}}{(n1)!}\left(\underset{m1}{}\frac{A_m}{m!}(\overline{\psi })^{m1}\right)_k𝐭,𝐭,\mathrm{},𝐭;𝐜(F_{g,n,d})_{g,n,d}.$$
Codim-2:
(7.4.1.2)
$$\frac{1}{2}\underset{g,n,d}{}\frac{Q^d\mathrm{}^{g1}}{n!}𝐭,\mathrm{},𝐭;\left[(f\iota )_{}\underset{m2}{}\frac{1}{m!}r_{node}^2ev_{node}^{}A_m\frac{\overline{\psi }_+^{m1}+(1)^m\overline{\psi }_{}^{m1}}{\overline{\psi }_++\overline{\psi }_{}}\right]_k𝐜(F_{g,n,d})_{g,n,d}.$$
Codim-0:
$$\begin{array}{cc}& \underset{g,n,d}{}\frac{Q^d\mathrm{}^{g1}}{n!}𝐭,\mathrm{},𝐭;(f_{}(ch(ev^{}F)Td^{}(L_{n+1})))_k𝐜(F_{g,n,d})_{g,n,d}\hfill \\ & =\underset{g,n,d}{}\frac{Q^d\mathrm{}^{g1}}{n!}f^{}𝐭,\mathrm{},f^{}𝐭,(ch(ev^{}F)Td^{}(L_{n+1}))_{k+1};𝐜(F_{g,n+1,d})_{g,n+1,d}^{}\hfill \\ & =\underset{g,n,d}{}\frac{Q^d\mathrm{}^{g1}}{n!}𝐭,\mathrm{},𝐭,(ch(F)Td^{}(L_{n+1}))_{k+1};𝐜(F_{g,n+1,d})_{g,n+1,d}^{}\hfill \\ & \underset{g,n,d}{}\frac{Q^d\mathrm{}^{g1}}{(n1)!}ch_{k+1}(F)_{orb}\left[\frac{𝐭(\overline{\psi })}{\overline{\psi }}\right]_+,𝐭,\mathrm{},𝐭;𝐜(F_{g,n,d})_{g,n,d},\hfill \end{array}$$
where we have used Lemma B.0.1. This is equal to the sum of the following four terms:
(7.4.1.3)
$$\underset{g,n,d}{}\frac{Q^d\mathrm{}^{g1}}{(n1)!}𝐭,\mathrm{},𝐭,(ch(F)Td^{}(L))_{k+1};𝐜(F_{g,n,d})_{g,n,d}^{}$$
(7.4.1.4)
$$\underset{g,n,d}{}\frac{Q^d\mathrm{}^{g1}}{(n1)!}ch_{k+1}(F)_{orb}\left[\frac{𝐭(\overline{\psi })}{\overline{\psi }}\right]_+,𝐭,\mathrm{},𝐭;𝐜(F_{g,n,d})_{g,n,d}$$
(7.4.1.5)
$$\frac{1}{2\mathrm{}}𝐭,𝐭,(ch(F)Td^{}(L))_{k+1};𝐜(F_{0,3,0})_{0,3,0}^{}$$
(7.4.1.6)
$$(ch(F)Td^{}(L))_{k+1};𝐜(F_{1,1,0})_{1,1,0}^{}.$$
Here $`\mathrm{}_{\mathrm{}}^{}`$ denotes invariants defined from moduli spaces of maps with the last marked point untwisted, and we use the property
$$f^{}𝐭(\overline{\psi }_j)=𝐭(\overline{\psi }_j)s_j\left[\frac{𝐭(\overline{\psi }_j)}{\overline{\psi }_j}\right]_+.$$
Since $`\overline{\psi }_j`$ are pulled back from $`\overline{}_{g,n}(X,d)`$, this follows from the case of schemes (see for instance ).
Observe that, by Lemma 2.3.8,
$$ch_{k+1}(F)_{orb}\left[\frac{𝐭(\overline{\psi })}{\overline{\psi }}\right]_+=ch_{k+1}(q^{}F)\left[\frac{𝐭(\overline{\psi })}{\overline{\psi }}\right]_+=ch_{k+1}(q^{}F)\left(\frac{𝐭(\overline{\psi })t_0}{\overline{\psi }}\right).$$
On a component $`𝒳_i`$, we have
$$\begin{array}{cc}& \underset{m1}{}\frac{A_m}{m!}z^{m1}|_{𝒳_i}=\underset{m1}{}\underset{0lr_i1}{}\frac{ch(F_i^{(l)})B_m(l/r_i)}{m!}z^{m1}\hfill \\ & =\underset{0lr_i1}{}ch(F_i^{(l)})\left(\underset{m1}{}\frac{B_m(l/r_i)}{m!}z^{m1}\right)=\underset{0lr_i1}{}ch(F_i^{(l)})\left(\frac{e^{\frac{l}{r_i}z}}{e^z1}\frac{1}{z}\right),\hfill \end{array}$$
and $`ch_{k+1}(q^{}F)|_{𝒳_i}=_{0lr_i1}ch_{k+1}(F_i^{(l)})`$.
For each $`l`$ we have
$$\left(\frac{e^{\frac{l}{r_i}z}}{e^z1}\frac{1}{z}\right)𝐭(z)+\frac{𝐭(z)t_0}{z}=\left(\frac{e^{\frac{l}{r_i}z}}{e^z1}𝐭(z)\frac{t_0}{z}\right)=\left[\frac{e^{\frac{l}{r_i}z}}{e^z1}𝐭(z)\right]_+.$$
Hence
$$\left(\underset{m1}{}\frac{A_m}{m!}\overline{\psi }^{m1}\right)_k|_{𝒳_i}𝐭(\overline{\psi })+ch_{k+1}(q^{}F)|_{𝒳_i}\left[\frac{𝐭(\overline{\psi })}{\overline{\psi }}\right]_+=\underset{0lr_i1}{}\left[\left(\frac{ch(F_i^{(l)})e^{\frac{l}{r_i}\overline{\psi }}}{e^{\overline{\psi }}1}\right)_k𝐭(\overline{\psi })\right]_+.$$
Therefore the sum of (7.4.1.1) and (7.4.1.4) is
$$\underset{g,n,d}{}\frac{Q^d\mathrm{}^{g1}}{(n1)!}\left[\left(\frac{_lch(F_i^{(l)})e^{\frac{l}{r_i}\overline{\psi }}}{e^z1}\right)_k𝐭\right]_+,𝐭,\mathrm{},𝐭;𝐜(F_{g,n,d})_{g,n,d}.$$
Also, $`(ch(F)Td^{}(L))_{k+1}=[(ch(F)\frac{Td^{}(L)}{\psi })_k\psi ]_+=[(\frac{ch(F)}{e^\psi 1})_k1\psi ]_+`$. Hence the sum of (7.4.1.1), (7.4.1.3) and (7.4.1.4) is
$$\underset{g,n,d}{}\frac{Q^d\mathrm{}^{g1}}{(n1)!}\left[\left(\frac{_lch(F_i^{(l)})e^{\frac{l}{r_i}\overline{\psi }}}{e^z1}\right)_k(𝐭1\overline{\psi })\right]_+,𝐭,\mathrm{},𝐭;𝐜(F_{g,n,d})_{g,n,d}.$$
Adding (7.1.1.2) to this, we get
(7.4.1.7)
$$\underset{g,n,d}{}\frac{Q^d\mathrm{}^{g1}}{(n1)!}\left[\left(\left(\frac{_lch(F_i^{(l)})e^{\frac{l}{r_i}\overline{\psi }}}{e^z1}\right)_k+\frac{ch_k((q^{}F)^{inv})}{2}\right)𝐪(\overline{\psi })\right]_+,𝐭,\mathrm{},𝐭;𝐜(F_{g,n,d})_{g,n,d}.$$
The restriction to $`𝒳_i`$ of the operator
(7.4.1.8)
$$\underset{\stackrel{m+h=k+1}{m,h0}}{}\frac{(A_m)_hz^{m1}}{m!}+\frac{ch_k((q^{}F)^{inv})}{2}$$
is
$$\underset{\stackrel{m+h=k+1}{m,h0}}{}\frac{(A_m)_hz^{m1}}{m!}|_{𝒳_i}+\frac{ch_k((q^{}F)^{inv}|_{𝒳_i})}{2}$$
(7.4.1.9)
$$=\left(\frac{_lch(F_i^{(l)})e^{\frac{l}{r_i}z}}{e^z1}\right)_k+\frac{ch_k((q^{}F)^{inv}|_{𝒳_i})}{2}.$$
Note that the operator (7.4.1.8) is infinitesimally symplectic by Corollary 4.1.5.
By , Example 1.3.3.1 (see Appendix C), the quantization of the $`pq`$-terms of the quadratic Hamiltonian of (7.4.1.8) applied to $`𝒟_s`$ gives (7.4.1.7).
It is straightforward to check that the $`q^2`$-term of the Hamiltonian of the operator (7.4.1.8) only comes from $`(A_0)_{k+1}/z=ch_{k+1}(q^{}F)/z`$. Using statements from Remark 2.4.3 we can calculate (7.4.1.5) directly, the answer is
$$\frac{1}{2\mathrm{}}_{I𝒳}t_0t_0ch_{k+1}(q^{}F)𝐜((q^{}F)^{inv}).$$
Then by Appendix C the quantization of the $`q^2`$-term yields exactly (7.4.1.5).
Now we handle the codim-2 terms (7.4.1.2), following the approach of . Pulling back to $`\overline{}_+\times \overline{}_{}`$ and $`\overline{}_{g1,n+\{+,\},d}`$ and using Lemma B.0.1 and Proposition 7.4.1, we express (7.4.1.2) as
(7.4.1.10)
$$\begin{array}{c}\hfill \frac{\mathrm{}}{2}\underset{g,n,d}{}\underset{g_1+g_2=g}{}\underset{n_1+n_2=n}{}\underset{d_1+d_2=d}{}\frac{Q^{d_1+d_2}\mathrm{}^{g_11+g_21}}{n_1!n_2!}\underset{a,b,c}{}𝐭,\mathrm{},𝐭,\frac{𝒪_{a,b,c}^{}\overline{\psi }_+^a}{\sqrt{𝐜((q^{}F)^{inv})}};𝐜(F_{g_1,n_1+1,d_1})_{g_1,n_1+1,d_1}\\ \hfill \times \frac{𝒪_{a,b,c,}^{\prime \prime }\overline{\psi }_{}^b}{\sqrt{𝐜((q^{}F)^{inv})}},𝐭,\mathrm{},𝐭;𝐜(F_{g_2,n_2+1,d_2})_{g_2,n_2+1,d_2}\\ \hfill +\frac{\mathrm{}}{2}\underset{g,n,d}{}\frac{Q^d\mathrm{}^{g11}}{n!}\underset{a,b,c}{}𝐭,\mathrm{},𝐭,\frac{𝒪_{a,b,c}^{}\overline{\psi }_+^a}{\sqrt{𝐜((q^{}F)^{inv})}},\frac{𝒪_{a,b,c,}^{\prime \prime }\overline{\psi }_{}^b}{\sqrt{𝐜((q^{}F)^{inv})}};𝐜(F_{g1,n+2,d})_{g1,n+2,d}.\end{array}$$
Here<sup>18</sup><sup>18</sup>18Note that the term $`_{m2}\frac{A_m}{m!}\frac{\overline{\psi }_+^{m1}+(1)^m\overline{\psi }_{}^{m1}}{\overline{\psi }_++\overline{\psi }_{}}`$ belongs to $`End(H^{}(I𝒳))[[\overline{\psi }_+,\overline{\psi }_{}]]`$, which is identified with $`H^{}(I𝒳)[[\overline{\psi }_+]]H^{}(I𝒳)[[\overline{\psi }_{}]]`$ using the pairing on $`H^{}(I𝒳)`$.
$$\underset{a,b}{}𝒪_{a,b}\overline{\psi }_+^a\overline{\psi }_{}^b=\left(\underset{m2}{}\frac{A_m}{m!}\frac{\overline{\psi }_+^{m1}+(1)^m\overline{\psi }_{}^{m1}}{\overline{\psi }_++\overline{\psi }_{}}\right)_{k1}(g^{\alpha \beta }\varphi _\alpha \varphi _\beta )H^{}(I𝒳)[[\overline{\psi }_+]]H^{}(I𝒳)[[\overline{\psi }_{}]],$$
$`g^{\alpha \beta }`$ is the matrix entry of the inverse of the matrix $`(g_{\alpha \beta })`$ with $`g_{\alpha \beta }=(\varphi _\alpha ,\varphi _\beta )_{orb}`$, and we write $`𝒪_{a,b}H^{}(I𝒳)H^{}(I𝒳)`$ in its Künneth decomposition:
$$𝒪_{a,b}=\underset{c}{}𝒪_{a,b,c}^{}𝒪_{a,b,c}^{\prime \prime },𝒪_{a,b,c}^{},𝒪_{a,b,c}^{\prime \prime }H^{}(I𝒳).$$
Due to twisted dilaton shift, we have
$$\frac{}{q_k^\alpha }=\frac{1}{\sqrt{𝐜((q^{}F)^{inv})}}\frac{}{t_k^\alpha }.$$
Comparing this with , Example 1.3.3.1 (see Appendix C), we find that (7.4.1.10) coincides with the quantization of the $`p^2`$-terms of the Hamiltonian of $`_{m,h0,m+h=k+1}\frac{(A_m)_h}{m!}z^{m1}+ch_k((q^{}F)^{inv})/2`$ applied to $`𝒟_s`$ (note that the Hamiltonian of $`(A_0)_{k+1}/z+(A_1)_k+ch_k((q^{}F)^{inv})/2`$ has no $`p^2`$-terms).
Putting the above together, we just proved
$$\frac{𝒟_s}{s_k}=\left[\left(\underset{m+h=k+1;m,h0}{}\frac{(A_m)_hz^{m1}}{m!}+\frac{ch_k((q^{}F)^{inv})}{2}\right)^{}+(\text{7.4.1.6})\right]𝒟_s.$$
Note that (7.4.1.6) is equal to $`C_k`$ defined in Section 7.1. This concludes the proof of (7.1.1.4).
## Appendix A A Grothendieck-Riemann-Roch formula for Stacks
Let $`𝒳`$ and $`𝒴`$ be Deligne-Mumford stacks with quasi-projective coarse moduli spaces. Let $`f:𝒳𝒴`$ be a proper morphism of Deligne-Mumford stacks. Assume that $`f`$ factors as
(A.0.0.1)
$$f=gi,$$
where $`i:𝒳𝒫`$ is a closed regular immersion and $`g:𝒫𝒴`$ is a smooth morphism (not necessarily representable). Define
$$T_f:=[i^{}T_{𝒫/𝒴}][N_{𝒳/𝒫}]K^0(𝒳).$$
It is easy to show that $`T_f`$ is independent of the factorization $`f=gi`$. There is a Grothendieck-Riemann-Roch formula for this kind of morphism, which is due to Toen . We begin with some definitions.
###### Definition A.0.1 ().
Define a map $`\rho :K^0(I𝒳)K^0(I𝒳)`$ as follows: If a bundle $`F`$ on $`I𝒳`$ is decomposed into a direct sum $`_\zeta F^{(\zeta )}`$ of eigenbundles $`F^{(\zeta )}`$ with eigenvalue $`\zeta `$, then
$$\rho (F):=\underset{\zeta }{}\zeta F^{(\zeta )}K^0(I𝒳).$$
###### Definition A.0.2 ().
Define $`\stackrel{~}{ch}:K^0(𝒳)H^{}(I𝒳)`$ to be the composite
$$K^0(𝒳)\stackrel{q_𝒳^{}}{}K^0(I𝒳)\stackrel{𝜌}{}K^0(I𝒳)\stackrel{ch}{}H^{}(I𝒳),$$
where $`q_𝒳:I𝒳𝒳`$ is the projection and $`ch`$ is the usual Chern character.
###### Definition A.0.3.
Define an operation $`\lambda _1`$ in K-theory as follows: for a vector bundle $`V`$, define $`\lambda _1(V):=_{a0}(1)^a\mathrm{\Lambda }^aV`$.
###### Definition A.0.4 (Todd class).
Define $`\stackrel{~}{Td}:K^0(𝒳)H^{}(I𝒳)`$ as follows: For a vector bundle $`E`$ on $`𝒳`$, $`q_𝒳^{}E`$ is decomposed into a direct sum $`(q_𝒳^{}E)^{inv}(q_𝒳^{}E)^{mov}`$ where $`(q_𝒳^{}E)^{inv}`$, the invariant part, is the eigenbundle with eigenvalue $`1`$, and $`(q_𝒳^{}E)^{mov}`$, the moving part, is the direct sum of eigenbundles with eigenvalues not equal to $`1`$. Define
$$\stackrel{~}{Td}(E):=\frac{Td((q_𝒳^{}E)^{inv})}{ch(\rho \lambda _1(((q_𝒳^{}E)^{mov})^{}))}.$$
The map $`\stackrel{~}{Td}`$ satisfies
$$\stackrel{~}{Td}(V_1+V_2)=\stackrel{~}{Td}(V_1)\stackrel{~}{Td}(V_2),\stackrel{~}{Td}(V_1V_2)=\frac{\stackrel{~}{Td}(V_1)}{\stackrel{~}{Td}(V_2)}.$$
Recall that a stack has the resolution property if every coherent sheaf is a quotient of a vector bundle (see for instance ).
###### Theorem A.0.5 (Grothendieck-Riemann-Roch formula ).
Let $`𝒳`$ and $`𝒴`$ be smooth Deligne-Mumford stacks with quasi-projective coarse moduli spaces and $`f:𝒳𝒴`$ a proper morphism which factors as (A.0.0.1). Assume that $`𝒳`$ and $`𝒴`$ have the resolution property. Let $`EK^0(𝒳)`$, then
$$\stackrel{~}{ch}(f_{}E)=If_{}(\stackrel{~}{ch}(E)\stackrel{~}{Td}(T_f)),$$
where $`f_{}`$ is the K-theoretic pushforward and $`If:I𝒳I𝒴`$ is the map induced by $`f`$.
###### Remark A.0.6.
The cohomological pushforward $`If_{}`$ of a non-representable morphism is defined by passing to a finite scheme cover of $`I𝒳`$, see .
Restricting to the distinguished component $`𝒴I𝒴`$, we obtain
###### Corollary A.0.7.
$$ch(f_{}E)=If_{}(\stackrel{~}{ch}(E)\stackrel{~}{Td}(T_f)|_{If^1(𝒴)}).$$
## Appendix B Properties of Virtual Bundles
In this appendix we discuss some properties of the virtual bundle $`F_{g,n,d}`$. First note that the fact that $`F_{g,n,d}`$ is well-defined can be seen by factoring $`f`$ as in (A.0.0.1) (which follows from the construction of the universal family in ). Note that $`f`$ is perfect, and resolution property implies that the $`K`$-theory of vector bundles coincides with the $`K`$-theory of perfect complexes.
We study how $`F_{g,n,d}`$ behaves under pulling back by the maps $`f:\overline{}_{g,n+1}(𝒳,d)^{}\overline{}_{g,n}(𝒳,d)`$, $`s_j:𝒟_j\overline{}_{g,n+1}(𝒳,d)^{}`$, and $`i:𝒵\overline{}_{g,n+1}(𝒳,d)^{}`$. Let $`\iota _{red}:\stackrel{~}{𝒵}^{red}\overline{}_{g,n+1}(𝒳,d)^{}`$ be the composition of double covering of $`𝒵^{red}`$ and the inclusion into $`\overline{}_{g,n+1}(𝒳,d)^{}`$. Similarly we can define $`\iota _{irr}:\stackrel{~}{𝒵}^{irr}\overline{}_{g,n+1}(𝒳,d)`$. By the definition of $`𝒵`$ it is the universal gerbe at node over $`f(𝒵)\overline{}_{g,n}(𝒳,d)`$. According to , Proposition 5.2.1, we have
$$\stackrel{~}{𝒵}^{red}=\underset{g_++g_{}=g,n_++n_{}=n,d_++d_{}=d}{}\overline{}_{g_+,n_++1}(𝒳,d_+)\times _{I𝒳}\overline{}_{g_{},n_{}+1}(𝒳,d_{}),$$
and
$$\stackrel{~}{𝒵}^{irr}=\overline{}_{g1,n+2}(𝒳,d)\times _{I𝒳\times I𝒳}I𝒳.$$
Therefore we may view $`\stackrel{~}{𝒵}^{red}`$ as the moduli stack which parametrizes pairs
(B.0.0.1)
$$(f_+:(𝒞_+,\{\mathrm{\Sigma }_i\}_{1in+}\{\mathrm{\Sigma }_+\})𝒳,f_{}:(𝒞_{},\{\mathrm{\Sigma }_i\}_{1in_{}}\{\mathrm{\Sigma }_{}\})𝒳),$$
where $`[f_\pm ]\overline{}_{g_\pm ,n_\pm }(𝒳,d_\pm )`$, such that
$$[f_+|_{\mathrm{\Sigma }_+}]=I([f_{}|_\mathrm{\Sigma }_{}])I𝒳.$$
Here $`I:I𝒳I𝒳`$ is the involution defined in Section 2.1
Similarly we may view $`\stackrel{~}{𝒵}^{irr}`$ as the moduli stack which parametrizes maps
(B.0.0.2)
$$[f:(𝒞,\{\mathrm{\Sigma }_i\}_{1in}\{\mathrm{\Sigma }_+,\mathrm{\Sigma }_{}\})𝒳]\overline{}_{g1,n+2}(𝒳,d),$$
such that
$$[f|_{\mathrm{\Sigma }_+}]=I([f|_\mathrm{\Sigma }_{}])I𝒳.$$
Let $`ev_{node}:\stackrel{~}{𝒵}I𝒳`$ denote the evaluation map at the marked point of gluing in the description of $`𝒵`$ above. More precisely, $`ev_{node}`$ is defined to map (B.0.0.1) to $`[f_+|_{\mathrm{\Sigma }_+}]I𝒳`$ and map (B.0.0.2) to $`[f|_{\mathrm{\Sigma }_+}]I𝒳`$.
###### Lemma B.0.1.
1. $`f^{}F_{g,n,d}=F_{g,n+1,d}|_{\overline{}_{g,n+1}(𝒳,d)^{}}.`$
2. $`\iota _{red}^{}F_{g,n+1,d}=p_+^{}F_{g_+,n_++1,d_+}+p_{}^{}F_{g_{},n_{}+1,d_{}}ev_{node}^{}(q^{}F)^{inv}.`$
3. $`\iota _{irr}^{}F_{g,n+1,d}=F_{g1,n+2,d}ev_{node}^{}(q^{}F)^{inv}`$.
###### Proof.
The proofs are similar to those of the corresponding statements in , . Let $`𝒳=[M/G]`$ be as in Assumption 2.5.9, where $`M`$ is a smooth quasi-projective variety and $`G`$ is a linear algebraic group. Choose a $`G`$-equivariant ample line bundle $`L`$ on $`M`$. The bundle $`F`$ corresponds to an equivariant vector bundle which we also denote by $`F`$. For $`N`$ sufficiently large we have the following exact sequence
$$0KerH^0(M,FL^N)FL^N0.$$
Tensoring with $`L^N`$ yields an exact sequence
$$0KerL^NH^0(M,FL^N)L^NF0.$$
Let $`A=H^0(M,FL^N)L^N`$ and $`B=KerL^N`$. These two bundles induce two vector bundles on $`𝒳`$ which we denote by $`𝒜`$ and $``$ respectively. The above exact sequence implies that $`F_{g,n,d}=𝒜_{g,n,d}_{g,n,d}`$.
If $`d0`$, then $`R^0f_{}ev_{n+1}^{}𝒜`$ and $`R^0f_{}ev_{n+1}^{}`$ both vanish for $`N`$ sufficiently large, and $`𝒜_{g,n,d},_{g,n,d}`$ are vector bundles.
We verify (2) for $`𝒜_{g,n,d}`$. Let $`T`$ be a scheme. Let
$$((f_+:𝒞_+𝒳),(f_{}:𝒞_{}𝒳))$$
be a $`T`$-valued point of
$$\overline{}_{g_+,n_++1}(𝒳,d_+)\times _{I𝒳}\overline{}_{g_{},n_{}+1}(𝒳,d_{}),$$
and $`f:𝒞𝒳`$ the stable map obtained by gluing. Denote by $`𝔱:𝒞T`$, $`𝔱_\pm :𝒞_\pm T`$ the structure maps, by $`\nu :𝒞_+𝒞_{}𝒞`$ the gluing morphism, and by $`\mathrm{\Theta }_{node}𝒞`$ the locus of the node formed by gluing. The restriction of $`\iota _{red}^{}𝒜_{g,n+1,d}`$ to the $`T`$-valued point $`(f:𝒞𝒳)`$ is
$$R^1𝔱_{}f^{}𝒜(R^0𝔱_{}(f^{}𝒜^{}\omega _𝒞))^{}.$$
The restriction of $`p_\pm ^{}𝒜_{g_\pm ,n_\pm +1,d_\pm }`$ to the $`T`$-valued point $`(f_\pm :𝒞_\pm 𝒳)`$ is
$$R^1𝔱_\pm f_\pm ^{}𝒜(R^0𝔱_\pm (f_\pm ^{}𝒜^{}\omega _{𝒞_\pm }))^{}.$$
The relative dualizing sheaves of $`𝒞,𝒞_+,𝒞_{}`$ are easily seen to fit into the following exact sequence,
$$0\omega _{𝒞/S}\nu _{}(\omega _{𝒞_+/S}\omega _{𝒞_{}/S})𝒪_{\mathrm{\Theta }_{node}}0.$$
Tensoring by $`f^{}𝒜^{}`$ and applying $`𝔱_{}`$ give the following exact sequence:
$$0R^0𝔱_{}(f^{}𝒜^{}\omega _{𝒞/T})R^0𝔱_+(f_+^{}𝒜^{}\omega _{𝒞_+/T})R^0𝔱_{}(f_{}^{}𝒜^{}\omega _{𝒞_{}/T})R^0𝔱_{}(f^{}𝒜^{}𝒪_{\mathrm{\Theta }_{node}})0.$$
Note that $`R^0𝔱_{}(f^{}𝒜^{}𝒪_{\mathrm{\Theta }_{node}})`$ is the sheave of sections of $`f^{}𝒜^{}`$ which are invariant under the action of the stabilizer group of the node. Therefore $`R^0𝔱_{}(f^{}𝒜^{}𝒪_{\mathrm{\Theta }_{node}})`$ is the restriction to the $`T`$-valued point $`((f_+:𝒞_+𝒳),(f_{}:𝒞_{}𝒳))`$ of $`ev_{node}^{}((q^{}𝒜^{})^{inv})`$ (c.f. the proof of Lemma 7.3.6). Dualizing this sequence then proves (2) for $`𝒜_{g,n,d}`$. We can prove it for $`_{g,n,d}`$ in the same way. (2) thus hold for $`F_{g,n,d}`$ since $`F_{g,n,d}=𝒜_{g,n,d}_{g,n,d}`$.
If $`d=0`$, then $`R^0f_{}ev_{n+1}^{}F`$ is a trivial bundle and $`R^1f_{}ev_{n+1}^{}F`$ is a vector bundle. The same argument can be applied to this case.
(1) and (3) can be proved by a similar approach, we omit the details. ∎
## Appendix C An Example of Quantized Operator
In this Appendix we reproduce the calculation in , Example 1.3.3.1.
Let $`A=Bz^m`$ be an infinitesimal symplectic transformation of $``$. Here $`B:H^{}(I𝒳)H^{}(I𝒳)`$ is a linear transformation. We write $`B`$ as a matrix $`(B_\beta ^\alpha )`$ using the basis $`\{\varphi _\alpha \}`$ of $`H^{}(I𝒳)`$. Put $`g_{\alpha \beta }=(\varphi _\alpha ,\varphi _\beta )_{orb}`$ and let $`g^{\alpha \beta }`$ denotes the matrix entry of the matrix inverse to $`(g_{\alpha \beta })`$. Then define $`B_{\alpha \beta }=g_{\alpha \gamma }B_\beta ^\gamma `$ and $`B^{\alpha \beta }=B_\gamma ^\alpha g^{\gamma \beta }`$.
A direct calculation shows that,
$$\widehat{A}=\frac{1}{2\mathrm{}}\underset{0km1}{}(1)^{k+m}B_{\alpha \beta }q_k^\beta q_{1km}^\alpha \underset{km}{}B_\beta ^\alpha q_k^\beta \frac{}{q_{k+m}^\alpha },\text{if }m<0,$$
and
$$\widehat{A}=\underset{k0}{}B_\beta ^\alpha q_k^\beta \frac{}{q_{k+m}^\alpha }+\frac{\mathrm{}}{2}\underset{0km1}{}(1)^kB^{\alpha \beta }\frac{}{q_k^\beta }\frac{}{q_{m1k}^\alpha },\text{if }m>0.$$
For $`m=0`$ we have $`\widehat{A}=_{k0}B_\beta ^\alpha q_k^\beta \frac{}{q_k^\alpha }`$.
We calculate
$$\begin{array}{cc}& \left(\underset{k}{}B_\beta ^\alpha q_k^\beta \frac{}{q_{k+m}^\alpha }\right)𝐪=\left(\underset{k}{}B_\beta ^\alpha q_k^\beta \frac{}{q_{k+m}^\alpha }\right)\left(\underset{l}{}q_l^\gamma \varphi _\gamma z^l\right)\hfill \\ & =\left[\underset{k}{}B_\beta ^\alpha q_k^\beta \varphi _\alpha z^{k+m}\right]_+=[A𝐪]_+.\hfill \end{array}$$
This explains the appearance of (7.4.1.7).
Now suppose $`m>0`$. We want to explain the double derivative terms in $`\widehat{A}`$ above, following , Example 1.3.3.1. Observe that the double derivative
$$\frac{}{q_k^\beta }\frac{}{q_{m1k}^\alpha }$$
is the bivector field corresponding to
$$\varphi _\beta \overline{\psi }_+^k\varphi _\alpha \overline{\psi }_{}^{m1k}_+_+(\text{identifying }\overline{\psi }_+,\overline{\psi }_{}\text{ with }z).$$
Note that for $`m1`$, we have
$$\underset{0km1}{}(1)^k\overline{\psi }_+^k\overline{\psi }_{}^{m1k}=\frac{\overline{\psi }_+^m+(1)^{m1}\overline{\psi }_{}^m}{\overline{\psi }_++\overline{\psi }_{}}.$$
Thus the term $`_{0km1}(1)^kB^{\alpha \beta }\frac{}{q_k^\beta }\frac{}{q_{m1k}^\alpha }`$ can be interpreted as the bivector field corresponding to
$$\frac{B\overline{\psi }_+^m+(1)^{m1}B\overline{\psi }_{}^m}{\overline{\psi }_++\overline{\psi }_{}}.$$
This explains the appearance of (7.4.1.2).
## Appendix D Cocycle calculation
In this appendix we calculate the cocycle (7.2.2.1). We begin with a lemma.
###### Lemma D.0.1.
Let $`𝒴`$ be a smooth proper Deligne-Mumford stack. Denote by $`q:I𝒴𝒴`$ the natural projection. Let $`A:H^{}(𝒴,)H^{}(𝒴,)`$ be a linear operator defined by a class $`aH^{}(𝒴,)`$, i.e. $`A(\gamma )=a\gamma `$. Then
$$\text{str}(A)=_{I𝒴}q^{}(a)e(T_{I𝒴}).$$
###### Proof.
Write the class $`a`$ as a sum of its degree zero part and positive degree part: $`a=a_01+a^{}`$ where $`a^{}H^{>0}(𝒴,)`$. Since $`H^{}(𝒴,)`$ is a graded ring, the operator of multiplication by a positive degree element of $`H^{}(𝒴,)`$ has super-trace $`0`$. So $`str(A)=str(a_01)=str(a_0id)`$.
We find that
$$\begin{array}{cc}\hfill str(id)& =\chi (I𝒴)\text{by the Lefschetz trace formula (see e.g. }\text{[9]}\text{)}\hfill \\ & =_{I𝒴}e(T_{I𝒴})\text{by Gauss-Bonnet (see e.g. }\text{[63]}\text{, Corollaire 3.44)}.\hfill \end{array}$$
Since $`q^{}a^{}e(T_{I𝒴})=0`$, we have
$$_{I𝒴}q^{}(a)e(T_{I𝒴})=_{I𝒴}q^{}(a_01)e(T_{I𝒴})=str(a_0id).$$
(7.2.2.1) is obtained by applying this Lemma to each component $`𝒳_i`$ of $`I𝒳`$, and use the definition of double inertia stack $`II𝒳:=I(I𝒳)=_iI𝒳_i`$. We denote the projection by $`Iq:II𝒳I𝒳`$.
## Appendix E Proof of (TRR)
In this Appendix we give a proof of the topological recursion relations (TRR) in genus $`0`$. In this proof we will use the moduli stack $`𝒦_{g,n}(𝒳,d)`$ instead of $`\overline{}_{g,n}(𝒳,d)`$. This is because the proof involves splitting nodal twisted curves along a node, and it is easier to express this using the stack $`𝒦_{g,n}(𝒳,d)`$. As pointed out in , Section 6.1.3, orbifold Gromov-Witten invariants defined using $`𝒦_{g,n}(𝒳,d)`$ agree with those defined using $`\overline{}_{g,n}(𝒳,d)`$. We refer to for properties of $`𝒦_{g,n}(𝒳,d)`$ used here.
Our proof is adopted from , Section VI.6.6. Let $`𝔐_{0,3+k}^{tw}`$ be the (Artin) stack of twisted curves of genus $`0`$ with $`3+k`$ marked gerbes (not trivialized) and denote by $`p:𝒦_{0,3+k}(𝒳,d)𝔐_{0,3+k}^{tw}`$ the forgetful morphism. For each partition $`\{4,5,\mathrm{},3+k\}=AB`$ with $`A,B`$ nonempty we consider the stack $`𝔇^{tw}(0;\{1\}A|0;\{2,3\}B)`$ defined in , Section 5.1. Put
$$𝔇^{tw}:=\underset{A,B;A{\scriptscriptstyle B}=\{4,5,\mathrm{},3+k\}}{}𝔇^{tw}(0;\{1\}A|0;\{2,3\}B).$$
There is a natural gluing map $`gl:𝔇^{tw}𝔐_{0,3+k}^{tw}`$ as defined in , Proposition 5.1.3. Form the following cartesian diagram
$$\begin{array}{ccc}𝒟(𝒳)& \stackrel{\mu }{}& 𝒦_{0,3+k}(𝒳,d)\\ & & p& & \\ 𝔇^{tw}& \stackrel{gl}{}& 𝔐_{0,3+k}^{tw}.\end{array}$$
Let $`𝒟^{tw}𝔐_{0,3+k}^{tw}`$ denote the image of $`𝔇^{tw}`$ under the map $`gl`$. Consider the forgetful maps $`𝔐_{0,3+k}^{tw}𝔐_{0,3+k}\overline{}_{0,3}`$, where the first map takes a twisted curve to its coarse curve, and the second map forgets all but the first three marked points and stabilizes the curves. Let $`L_1`$ be the line bundle over $`𝔐_{0,3+k}^{tw}`$ obtained by pulling back the first universal cotangent line bundle over $`𝔐_{0,3+k}`$, and $`L_1^{}`$ the line bundle over $`𝔐_{0,3+k}^{tw}`$ obtained by pulling back the first universal cotangent line bundle over $`\overline{}_{0,3}`$. (We slightly abuse notations here.) It is not hard to see that there is an exact sequence
$$0L_1^{}L_1𝒪_{𝒟^{tw}}0.$$
A standard intersection theory result (see e.g. , Chapter VI, equation (6.19)) shows that for any cycle class $`\alpha `$ on $`𝒦_{0,3+k}(𝒳,d)`$ we have
$$c_1(p^{}L_1)\alpha =c_1(p^{}L_1^{})\alpha +\mu _{}gl^!\alpha .$$
Take $`\alpha =[𝒦_{0,3+k}(𝒳,d)]^{vir}`$ and use the fact that $`c_1(L_1^{})=0`$ (because $`\overline{}_{0,3}`$ is a point), we get
(E.0.0.1)
$$\overline{\psi }_1[𝒦_{0,3+k}(𝒳,d)]^{vir}=\mu _{}(gl^![𝒦_{0,3+k}(𝒳,d)]^{vir}).$$
According to , Proposition 5.2.2, we have
$$𝔇^{tw}(0;\{1\}A|0;\{2,3\}B)\times _{𝔐_{0,3+k}^{tw}}𝒦_{0,3+k}(𝒳,d)\underset{d_1+d_2=d}{}𝒦_{0,\{1\}A}(𝒳,d_1)\times _{\overline{I}𝒳}𝒦_{0,\{2,3\}B}(𝒳,d_2),$$
where $`\overline{I}𝒳`$ is the rigidified inertia stack of $`𝒳`$ (see , Section 3.4). The diagonal map $`\delta :\overline{I}𝒳\overline{I}𝒳\times \overline{I}𝒳`$ fits into the following cartesian diagram
$$\begin{array}{ccc}𝒦_{0,\{1\}A}(𝒳,d_1)\times _{\overline{I}𝒳}𝒦_{0,\{2,3\}B}(𝒳,d_2)& & 𝒦_{0,\{1\}A}(𝒳,d_1)\times 𝒦_{0,\{2,3\}B}(𝒳,d_2)\\ ev_{node}& & ev_{}\times \mathrm{ev}_{}& & \\ \overline{I}𝒳& \stackrel{\delta }{}& \overline{I}𝒳\times \overline{I}𝒳.\end{array}$$
By the splitting result (, Proposition 5.3.1), we get
(E.0.0.2)
$$gl^![𝒦_{0,3+k}(𝒳,d)]^{vir}=\underset{A{\scriptscriptstyle B}=\{4,5,\mathrm{},3+k\};d_1+d_2=d}{}\delta ^!([𝒦_{0,\{1\}A}(𝒳,d_1)]^{vir}\times [𝒦_{0,\{2,3\}B}(𝒳,d_2)]^{vir}).$$
We may apply (E.0.0.2) to (E.0.0.1) and view the resulting equality in homology via cycle map. Integrate the resulting equality against $`\gamma =\varphi _{\alpha _1}\overline{\psi }_1^{k_1}_{i=2}^{3+k}\varphi _{\alpha _i}\overline{\psi }_i^{k_i}`$ and use an identification of $`H^{}(I𝒳)`$ and $`H^{}(\overline{I}𝒳)`$ (c.f. section 6.1.3), we get
$$\varphi _{\alpha _1}\overline{\psi }_1^{k_1+1}\underset{i=2}{\overset{3+k}{}}\varphi _{\alpha _i}\overline{\psi }_i^{k_i}_{0,3+k,d}$$
$$=\underset{A{\scriptscriptstyle B}=\{4,\mathrm{},3+k\},d=d_1+d_2}{}\underset{\alpha }{}\pm \varphi _{\alpha _1}\overline{\psi }_1^{k_1},\underset{iA}{}\varphi _{\alpha _i}\overline{\psi }_i^{k_i},\varphi _a_{0,|A|+2,d_1}\varphi ^a,\varphi _{\alpha _2}\overline{\psi }_2^{k_2},\varphi _{\alpha _3}\overline{\psi }_3^{k_3},\underset{iB}{}\varphi _{\alpha _i}\overline{\psi }_i^{k_i}_{0,|B|+3,d_2}.$$
Here the sign come from the possibly different ordering of odd cohomology classes between the left and right sides. (TRR) follows as this is the equality of coefficients of the corresponding terms on the left and right sides of (TRR). |
warning/0506/astro-ph0506470.html | ar5iv | text | # Evidence for Spectropolarimetric Diversity in Type Ia Supernovae
## 1 Introduction
The substantial homogeneity of Type Ia supernovae (SNe Ia), together with the fact that the peak luminosity of individual objects can be observationally determined through comparisons with well-calibrated samples of nearby SNe Ia, has propelled them to “gold standard” status among extragalactic distance indicators. Precise distance measurements out to $`z1.7`$ have been made (Riess et al. 2001, 2004), revealing the surprising cosmological result that the expansion rate of the universe is currently accelerating (Riess et al. 1998a; Perlmutter et al. 1999); see Filippenko (2004, 2005) for extensive reviews. Although the progenitor systems have not yet been conclusively identified, the general consensus is that SNe Ia arise from carbon-oxygen white dwarfs (CO WDs) that accrete matter through some mechanism until they achieve a density of $`3\times 10^9\mathrm{g}\mathrm{cm}^3`$ in their centers, leading to a runaway thermonuclear reaction that incinerates the star (Woosley & Weaver 1986). This occurs, coincidentally, when a CO WD’s mass is nearly the Chandrasekhar limit of $``$1.4 $`M_{}`$. Simulations of exploding CO WDs are able to reproduce the main spectral and photometric characteristics of SNe Ia (e.g., Leibundgut 2000, and references therein), lending theoretical support to this scenario.
However, many questions remain concerning the SN Ia progenitors and explosion mechanism. What is the nature of the “donor” that is responsible for providing the material that the WD accretes? Is there more than one channel by which the accretion can take place? Where does the thermonuclear runaway begin inside the WD, and how does it propagate throughout the star? As with so many things, the “devil is in the details,” and there is a growing sentiment that the answers to these fundamental questions may come from careful study of the differences seen among SNe Ia, rather than from the similarities alone.
Spectropolarimetry offers the only direct probe of early-time SN geometry, and thus is an important diagnostic tool for discriminating among SN Ia progenitor systems and theories of the explosion physics. The essential idea is this: A hot, young SN atmosphere is dominated by electron scattering, which by its nature is highly polarizing. For an unresolved source that has a spherical distribution of scattering electrons, the directional components of the electric vectors of the scattered photons cancel exactly, yielding zero net linear polarization. Any asymmetry in the distribution of the scattering electrons, or of absorbing material overlying the electron-scattering atmosphere, results in incomplete cancellation, and produces a net polarization (see, e.g., Leonard & Filippenko 2005).
Initial broad-band polarimetry studies found SNe Ia to possess zero or, at most, very weak ($`<0.2\%`$) intrinsic polarization (Wang et al. 1996), suggesting a high degree of symmetry for their scattering atmospheres. More recent spectropolarimetric studies capable of resolving individual line features, however, are revealing a complex picture, with both continuum and line polarization now convincingly established for at least a subset of the SN Ia population.
To date, three SNe Ia have been examined in detail with spectropolarimetry at early times: SN 1999by (Howell et al. 2001), SN 2001el (Kasen et al. 2003; Wang et al. 2003), and, most recently, SN 2004dt (Wang et al. 2005a). In this paper we present single-epoch spectropolarimetry of four SNe Ia: SN 1997dt, SN 2002bf, SN 2003du, and SN 2004dt, obtained about $`21`$, 3, 18, and 4 days (respectively) after maximum light.
A particular motivation for this multi-object investigation is to attempt to link spectral and photometric peculiarities of individual SNe Ia with their spectropolarimetric characteristics. Accordingly, the four events span a range of properties: SN 1997dt is likely somewhat subluminous, SN 2002bf and SN 2004dt exhibit unusually high-velocity absorption lines (in the case of SN 2002bf, the highest ever seen in an SN Ia for the epochs considered), and SN 2003du is slightly overluminous.
This paper is organized as follows. We briefly review the present photometric, spectroscopic, and spectropolarimetric state of knowledge of SNe Ia in § 2, focusing particular attention on objects sharing characteristics with the SNe Ia included in our spectropolarimetric survey. We describe and present the spectropolarimetry in § 3; to assist in the classification, we also include optical photometry and additional spectroscopy for two of the events, SN 2002bf and SN 2003du. We analyze the data in § 4, and present our conclusions in § 5. Note that preliminary discussions of the spectropolarimetry of SN 1997dt and SN 2002bf have been given by Leonard et al. (2000b) and Filippenko & Leonard (2004), respectively.
## 2 Background
### 2.1 Photometric Properties of SNe Ia
Type Ia SNe typically rise to peak $`B`$-band brightness in about 20 days, decline by about $`3`$ mag over the next $`35`$ days, and then settle into a nearly constant descent of $``$1.55 mag (100 day)<sup>-1</sup> for the next year. However, it is now confirmed beyond doubt that deviations exist from the central trend; see, for example, Phillips et al. (1999). Some events rise and fall more slowly, producing “broader” light curves (e.g., SN 1991T), whereas others rise and fall more quickly, yielding “narrower” light curves (e.g., SN 1991bg). The width of the light curves near peak correlates strongly with luminosity, in the sense that broader light curves generally indicate intrinsically brighter objects (but see Jha et al. 2005, for exceptions), with a total spread of about a factor of sixteen in absolute peak $`B`$-band brightness among the population (Altavilla et al. 2004).
The ability to individually determine the luminosity of SNe Ia through examination of their light-curve shapes has been the primary driver for cosmological applications. There are several photometric calibration techniques in common use that correlate the observed light-curve shape, or supernova color, with luminosity (see Wang et al. 2005b, and references therein). The simplest calibration technique involves measuring the decline in $`B`$ during the first 15 days after maximum $`B`$-band brightness, $`\mathrm{\Delta }m_{15}(B)`$. A typical value is $`\mathrm{\Delta }m_{15}(B)1.1`$ mag, which corresponds to $`M_B=19.26\pm 0.05`$ mag according to the calibrations of Hamuy et al. (1996). Underluminous, SN 1991bg-like objects yield values up to $`\mathrm{\Delta }m_{15}(B)=1.94`$ mag, while overluminous, SN 1991T-like objects go as low as $`\mathrm{\Delta }m_{15}(B)=0.81`$ mag (Altavilla et al. 2004). There is a continuum of values in between the two extremes.
For the two SNe with optical photometry presented in this paper, we shall use the “multicolor light-curve shape” method (MLCS; e.g., Riess et al. 1996), which has been revised (and hereafter referred to as MLCS2k2) by Jha (2002) and S. Jha et al. (in preparation) to include $`U`$-band light curves from Jha et al. (2005), a more self-consistent treatment of extinction, and an improved determination of the unreddened SN Ia color. This method is based on determining $`\mathrm{\Delta }`$ (a dimensionless number related to magnitude that parameterizes the light-curve shape), $`t_0`$ (the time of $`B`$-band maximum light), $`\mu _0`$ (the distance modulus for an assumed value of $`H_0`$, here taken to be $`65\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$), and $`A_V^0`$ (the visual extinction at $`t_0`$).
The mechanism that produces the dispersion in SN Ia luminosity is not known; present speculations range from different progenitor systems to differences in the explosion mechanism and flame propagation (e.g., deflagration, detonation, delayed-detonation, or, most recently, gravitationally confined detonation; see, respectively, Nomoto et al. 1984; Arnett 1969; Khokhlov 1991; Plewa et al. 2004). Some have further speculated that global asymmetries in the expanding ejecta may result in viewing-angle dependent luminosity (Wang et al. 2003). One observational fact that all theories must confront is that the most luminous SNe Ia have thus far only been seen in late-type galaxies; the faintest objects tend to prefer elliptical galaxies, but have also been found in spirals (e.g., Howell 2001; Benetti et al. 2005, and references therein).
### 2.2 Spectroscopic Properties of SNe Ia
Early-time spectra of SNe Ia (see Filippenko 1997 for a review) typically exhibit lines of intermediate-mass elements (IMEs), such as magnesium, silicon, sulfur, and calcium, with some contribution from iron-peak elements, especially at near-ultraviolet wavelengths. As time progresses, and the photosphere recedes deeper into the ejecta, lines of Fe come to dominate the spectrum. This spectral evolution suggests a burning front that incinerates some of the progenitor’s carbon and oxygen all the way to the iron peak deep inside the ejecta, but then leaves the outer layers only partially burned. It is interesting to note that overluminous SNe Ia show enhanced Fe features at early times, whereas subluminous ones show weak early-time Fe features.
At early times (e.g., $`20\mathrm{d}<t<20\mathrm{d}`$ from the date of maximum $`B`$ brightness, $`B_{\mathrm{max}}`$), the spectra of typical SNe Ia evolve rapidly and with such uniformity that it is possible to determine the age of an event relative to the date of $`B_{\mathrm{max}}`$ to within $``$2 days from a single spectrum alone (Riess et al. 1997; Foley et al. 2005). The spectroscopic peculiarities of subluminous and overluminous events currently preclude their “spectral feature ages” from being derived accurately through comparison with average SN Ia spectra. For instance, pre-maximum spectra of overluminous events lack the strong Si II $`\lambda 6355`$ absorption that is so prominent in normal and underluminous SN Ia spectra (Filippenko et al. 1992b).
Another point of distinction among SN Ia spectra comes from the blueshifts of the spectral lines. Increasing attention is being paid to a small but growing group of “high-velocity” (HV) SNe Ia, whose spectra around maximum light are characterized by unusually broad and highly blueshifted absorption troughs in many line features, indicating optically thick ejecta moving about 4000–5000 km s<sup>-1</sup> faster than is typically seen for SNe Ia (Branch 1987; Benetti et al. 2005). Well-studied examples of HV SNe Ia include SN 1983G (Branch & van den Bergh 1993, and references therein), SN 1984A (Branch 1987; Barbon et al. 1989), SN 1997bp (Anupama 1997), SN 1997bq (Lentz et al. 2001, and references therein), SN 2002bo (Benetti et al. 2004), SN 2002dj (Benetti et al. 2004), and SN 2004dt (Wang et al. 2005a). Wang et al. (2005a) note that in a spectrum of one HV SN Ia, SN 2004dt, some line features do not possess abnormally high velocity, such as those identified with S II. Since two of our objects, SN 2002bf and SN 2004dt, are HV SNe Ia, we briefly review the salient features that are known about this class of objects.
The most thorough and recent study of the diversity of SN Ia expansion velocities is that of Benetti et al. (2005), who apply a statistical treatment to data from 26 SNe Ia and find that HV SNe Ia indeed make up a kinematically distinct group. Its members have normal peak luminosity (based on their $`\mathrm{\Delta }m_{15}(B)`$ values: $`1.09\mathrm{\Delta }m_{15}(B)1.37`$ mag), reside in all types of host galaxies (e.g., both ellipticals and spirals, but with a preference for later types), and have very strong Si II features. They are further distinguished by a large temporal velocity gradient, $`\dot{v}`$ $`>70\mathrm{km}\mathrm{s}^1\mathrm{day}^1`$, where $`\dot{v}`$ is defined to be the average daily rate of decrease of the expansion velocity between maximum light and the time the Si II $`\lambda 6355`$ feature disappears. HV SNe Ia appear to have similar photometric characteristics to ordinary Type Ia events, although there is some indication of subtle differences in their color evolution (Benetti et al. 2004).
Whether HV SNe Ia represent the extreme end of a continuum of more typical SNe Ia, or require a different explosion mechanism, progenitor system, or explosion physics, remains controversial. What is certain is that any model must have significant optical depth in the IMEs at high velocities at early times. Branch & van den Bergh (1993) originally proposed that HV SNe Ia may simply result from more energetic explosions. However, Benetti et al. (2004) show that the Fe nebular lines have velocities comparable to those of normal SNe Ia, which tends to weaken the argument for higher overall kinetic energy. There is also evidence from early-time spectra of SN 2002bo that HV SNe Ia possess very little unburned carbon in their outer layers, in contrast with more typical SNe Ia like SN 1994D (Benetti et al. 2004). The high velocities of the IMEs coupled with the lack of primordial carbon may indicate that burning to Si penetrates to much higher layers in HV SNe Ia than it does in more normal events. This is a feature of some delayed-detonation models, such as those studied by Lentz et al. (2001). However, Wang et al. (2005a) argue that a strong detonation wave is unlikely to generate the clumpy and asymmetrically distributed silicon layer that is inferred from pre-maximum spectropolarimetric observations of SN 2004dt. As another alternative, Benetti et al. (2004) propose that IMEs produced at deeper layers may simply be more efficiently mixed outward in HV SNe Ia than is typical. A prediction of this mechanism is that while the IME products of explosive carbon burning will be mixed in with the primordial C and O in the outer layers, the C and O should also be mixed inward within these events, and exist at lower velocities than are normally seen.
It is appropriate at this point to mention the recently proposed gravitationally confined detonation (GCD) model of Plewa et al. (2004), which stands as an intriguing alternative to the standard deflagration or delayed-detonation scenarios. The mechanism involves the slightly off-center ignition of a deflagration that produces a buoyancy-driven bubble of material that reaches the stellar surface at supersonic speeds, where it laterally accelerates the outer stellar layers. This material, gravitationally confined to the white dwarf, races around the star and, in $``$2 s, converges at a point opposite to the location of the bubble’s breakout, creating conditions capable of igniting the nuclear fuel and triggering a detonation that can incinerate the white dwarf and result in an energetic explosion.
Of interest to the present study is the recent work of Kasen & Plewa (2005), who investigate the spectral, and spectropolarimetric, consequences of the GCD model. They focus their investigation on the interaction of the expanding ejecta with an ellipsoidal, metal-rich extended atmosphere formed from the bubble of deflagration products (taken to be 57% Si, 27% S, 7.1% Fe, and 2.7% Ca, plus smaller amounts of other metals; see Khokhlov et al. 1993), and find that a dense, optically thick pancake of metal-rich material is formed at potentially large velocity on the side of the ejecta where the bubble emerged. For low atmosphere masses (e.g., resulting from a bubble of mass $`0.008M_{}`$), the pancake of material spans the velocity range 17,000–28,000 km s<sup>-1</sup> and is geometrically detached from the bulk of the SN ejecta. This might explain the detached, high-velocity Ca II near-infrared (IR) triplet absorption seen in pre-maximum spectra of some SNe Ia (e.g., SN 2001el, see § 2.3.3). For larger bubble and, hence, atmosphere masses ($`m_{\mathrm{atm}}0.016M_{}`$), the absorbing pancake moves at lower velocities (e.g., 10,000–21,000 km s<sup>-1</sup> for $`m_{\mathrm{atm}}=0.08M_{}`$) and could blend with the region of IMEs in the SN ejecta and potentially increase the blueshift of several of the spectral features. This could provide an orientation-dependent explanation for the origin of HV SNe Ia, whereby sight lines in which the pancake more completely blocks the photosphere produce the anomalously large velocities. The spectropolarimetric consequences of this model are discussed in § 2.3.1.
### 2.3 Spectropolarimetric Properties of SN Ia
#### 2.3.1 Supernova Polarization Mechanisms
The first definitive proof that some SNe are polarized came from observations of the Type II SN 1987A (Cropper et al. 1988), which exhibited a modest temporal increase in continuum polarization during the photospheric phase, as well as sharp polarization modulations across strong P-Cygni flux lines (Jeffery 1991b). While intrinsic polarization has now been established in over a dozen SNe (for recent reviews, see Wheeler 2000; Filippenko & Leonard 2004; Leonard & Filippenko 2005), including at least three SNe Ia, the exact origin of both continuum and line polarization remains controversial.
There are essentially two mechanisms by which supernova continuum polarization is thought to be produced: (1) a globally aspherical photosphere and electron-scattering atmosphere (e.g., Shapiro & Sutherland 1982; Höflich 1991; Jeffery 1991a; Leonard et al. 2000a; Wang et al. 2001), and (2) ionization asymmetry produced by the decay of asymmetrically distributed radioactive $`{}_{}{}^{56}\mathrm{Ni}`$, perhaps flung out into (Chugai 1992; Höflich et al. 2001) or beyond (Kawabata et al. 2002; Leonard et al. 2002a) the expanding ejecta in clumps. The simplest, and most well-studied, globally aspherical geometry is that of an ellipsoid, and we shall make frequent reference to the “ellipsoidal model” in the following discussion. In the second model, in which an ionization asymmetry is present, continuum polarization is generated by light from the (either spherical or aspherical) photosphere scattering off of asymmetrically distributed free electrons that exist in clouds surrounding clumps of radioactive $`{}_{}{}^{56}\mathrm{Ni}`$. Variations on both of these polarization mechanisms are also possible. For instance, an aspherical distribution of $`{}_{}{}^{56}\mathrm{Ni}`$ could also result in SN polarization by providing an asymmetry in a source of optical photons (produced by the thermalization of $`\gamma `$-rays) relative to the scattering medium.
To explain polarization modulations seen across spectral lines, it is important to differentiate between the emission peaks and blueshifted absorption troughs that are characteristic of P-Cygni profiles in total-flux spectra of SNe. Line peaks have usually been assumed to consist of intrinsically unpolarized photons. This is because although resonance scattering by a line is an inherently polarizing process (Jeffery 1991a), directional information for scattered photons is lost in an SN atmosphere since the timescale for randomizing collisional redistribution of the relative level populations within the fine structure of the atomic levels of a line transition is much shorter than the characteristic timescale for absorption and reemission in a line (Höflich et al. 1996). The assumption of intrinsically unpolarized emission lines is also commonly used to derive the interstellar polarization (ISP; § 2.3.2). Note, though, that if ionization asymmetry exists above the photosphere (from, e.g., clumps of radioactive $`{}_{}{}^{56}\mathrm{Ni}`$), then even intrinsically unpolarized emission-line photons may become polarized by electron scattering within the SN atmosphere.
For SNe Ia, line-blanketing, due largely to Fe, is particularly severe at wavelengths below $``$5000 Å, where theoretical models (Howell et al. 2001) suggest that nearly complete depolarization of any “continuum” light may be assumed. Conversely, the broad spectral region $`6800\lambda 7800`$ Å is largely devoid of line opacity (Kasen et al. 2004; Howell et al. 2001), at least near maximum light, and thus may give a good indication of true continuum polarization level.
The explanation of polarization changes in absorption troughs is controversial. One well-studied hypothesis, which follows logically from the ellipsoidal model, is that selective blocking of more forward-scattered and, hence, less polarized, light in P-Cygni absorption troughs results in trough polarization increases (Jeffery 1991a; Leonard et al. 2000b, 2001), potentially producing “inverted” P-Cygni polarization profiles when coupled with emission peak depolarization. Some SN modelers have claimed that this “geometrical dilution” mechanism, however, is rather poor at polarizing SN Ia light, and that polarization decreases should actually be seen in absorption troughs, since much of the light reaching the observer in those spectral regions has been absorbed and reemitted by the line (Howell et al. 2001).
A key point is that, in either case, to produce line trough polarization changes under the ellipsoidal model, some continuum polarization must exist. In fact, Leonard & Filippenko (2001) show that the strength of a line trough polarization feature can be used to place a lower bound on the true intrinsic continuum polarization level under the ellipsoidal model according to
$$p_{cont}\frac{\mathrm{\Delta }p_{tot}}{(I_{cont}/I_{trough})1},$$
(1)
where $`I_{cont}`$ is the interpolated value of the continuum flux at the location of the line trough, $`I_{trough}`$ is the total flux at the line’s flux minimum, and $`\mathrm{\Delta }p_{tot}`$ is the total polarization change observed in the line trough. Thus, a critical test of the ellipsoidal model is to see whether the continuum polarization observed in the spectral region 6800–7800 Å is sufficient to explain observed line trough polarization changes. Lack of significant continuum polarization in an object with strong line features argues against the ellipsoidal model. An additional prediction is that no rotation of the polarization angle (PA) should be seen across the line, since the continuum and line-forming regions share the same geometry.
A final way to produce a polarization change in a line trough is through asymmetry in the distribution of elements in the ejecta material located above the photosphere along the line-of-sight (l-o-s), hereafter referred to as the “clumpy ejecta” model. Asymmetry in the distribution of material with significant optical depth may unevenly block the underlying photospheric light, thereby producing a net polarization change and/or PA rotation in a line trough, even when the photosphere is spherical. The GCD model discussed in § 2.2 provides a natural mechanism by which to generate line-trough polarization through this model, since it predicts an optically thick pancake of high-velocity, metal-rich material overlying the photosphere on the side in which the bubble emerged. In particular, if the GCD model provides the correct mechanism by which to produce HV SNe Ia, then polarization changes in the absorption troughs of the strong, HV metal lines in these objects should be particularly pronounced since the l-o-s necessarily intercepts a substantial fraction of the pancake.
For SNe Ia, an interesting alternative to the simple ellipsoidal model for producing both line and continuum polarization is presented by Kasen et al. (2004), who explore the polarization consequences of a conical hole in the ejecta due to the interaction with a companion star, which we shall refer to as the “ejecta-hole” model. By considering various viewing angles and hole sizes, Kasen et al. (2004) demonstrate that both continuum and line polarization and PA changes can be generated. In general, viewing angles almost directly down the hole yield low continuum polarization with polarization increases in strong line troughs, whereas sight lines more nearly perpendicular to the hole result in larger continuum polarization (prominently seen in the spectral region 6800–7800 Å) and “inverted P-Cygni” line polarization profiles; this latter case is nearly indistinguishable from the predictions of the simple ellipsoidal model.
Ionization asymmetry is also capable of generating both continuum polarization as well as PA and polarization level changes through line features, since both continuum and line photons will scatter off of concentrations of free electrons in (or beyond) the ejecta. An impressive example of this mechanism potentially being at work is given by Chugai (1992) for the case of SN 1987A, in which the spectropolarimetry data, as well as asymmetries in the flux line profiles, are convincingly reproduced by the effects of two clumps of $`{}_{}{}^{56}\mathrm{Ni}`$ in the far (receding) hemisphere of the ejecta.
#### 2.3.2 Removing Interstellar Polarization
A problem that plagues interpretation of all SN polarization measurements is proper removal of the ISP. Since directional extinction resulting from aspherical interstellar dust grains aligned by some mechanism along the l-o-s to an SN can contribute a large polarization to the observed signal, an attempt must be made to remove it prior to analyzing SN spectropolarimetry data. This is notoriously difficult, although a number of different techniques have been advanced over the years, which we here summarize.
An excellent way to derive Galactic ISP is through observations of distant, intrinsically unpolarized, “probe stars” close to (within $``$0.5$`\mathrm{°}`$) the l-o-s to the SN (e.g., Leonard et al. 2002a, b). Deriving the total ISP, which includes the contribution from dust in the host galaxy, however, is more difficult.
The most basic technique is to place limits on the ISP from reddening considerations. Since the same dust that polarizes starlight should redden it as well, it seems logical to expect a correlation between reddening and ISP. Because the alignment of dust grains is not total (or has multiple preferred orientations due to non-uniformity of the magnetic field along the l-o-s), and grains are probably only moderately elongated particles, it is not surprising that, through the analysis of thousands of reddened Galactic stars, only an upper bound on the polarization efficiency of Galactic dust can be derived (Serkowski et al. 1975): $`\mathrm{ISP}/E_{BV}<9.0\%\mathrm{mag}^1`$. However, the polarization efficiency of the dust in external galaxies is not well studied, and in one of the few investigations carried out to date, Leonard et al. (2002b) find compelling evidence for polarization efficiency well in excess of the empirical Galactic limit for dust in NGC 3184 along the l-o-s to the Type II-P SN 1999gi. It is not clear at this point whether meaningful constraints can thus be placed on the ISP of extragalactic SNe that are significantly reddened by host-galaxy dust. Nonetheless, total reddening arguments are still often used to set “reasonable” limits on the ISP.
Another method to get a handle on the ISP is to assume axisymmetry for any SN asphericity, a situation that reveals itself through a straight-line distribution of points when the spectropolarimetry is plotted in the $`q`$$`u`$ plane. If axisymmetry exists, the ISP is constrained to lie along the axis defined by the line (Howell et al. 2001). Its absolute value, however, is uncertain without additional input.
A more precise technique relies on the theoretical expectation that unblended emission lines consist of unpolarized light, and that any polarization observed in emission-line photons comes from the ISP (Jeffery 1991b; Tran et al. 1997; Wang et al. 1996; Leonard et al. 2001). This assumption is thought to be most valid at early times, when any $`{}_{}{}^{56}\mathrm{Ni}`$ concentrations are likely to be below the photosphere. Note, though, that to use this method, care must be taken to isolate the emission-line photons from the underlying continuum light (e.g., Tran et al. 1997), which may be intrinsically polarized.
A related method is to assume that a particular spectral region is intrinsically completely unpolarized, with all of the observed polarization therefore coming from ISP. Some empirical support exists that, in some objects, this may be the case for very strong emission lines (Kawabata et al. 2002; Leonard et al. 2002a; Wang et al. 2004). For SNe Ia in particular, it has sometimes been assumed that the far blue spectral region (e.g., typically $`<5000`$ Å) satisfies this criterion due to the effect of the heavy line-blanketing and multiple, depolarizing, line scatters, largely due to iron-group elements. This qualitative expectation is demonstrated quantitatively by Howell et al. (2001), who present model polarization spectra resulting from delayed-detonation models in realistic SN Ia atmospheres. More recently, Wang et al. (2005a) update this approach by choosing only specific spectral regions at blue wavelengths ($`\lambda <5000`$ Å) that are not dominated by obvious individual line features in either flux or polarization for the ISP determination, rather than the arbitrary blue edge of the spectrum, as adopted by Howell et al. (2001). The central idea here is that while strong individual line features might impart their own polarization signature (e.g., through geometrical dilution in ellipsoidal models, or through blocking of the parts of the photosphere in clumpy-ejecta models), the spectral regions in between specific features probably decrease any effective continuum polarization by the numerous overlapping spectral lines at blue wavelengths.
Finally, if polarimetry is obtained after the electron-scattering optical depth of the atmosphere has dropped well below unity, then it may be adequate to assume the observed polarization to be due entirely to ISP across the whole spectrum (e.g., Wang et al. 2003).
Improperly removed, ISP can increase or decrease the derived intrinsic polarization, and it can change “valleys” into “peaks” (or vice versa) in the polarization spectrum. Since it is so difficult to be certain of accurate removal of ISP, it is generally safest to focus on (a) temporal changes in the polarization with multiple-epoch data, (b) distinct line features in spectropolarimetry having high signal-to-noise ratio (S/N), and (c) continuum polarization unlike the characteristic “Serkowski-law” wavelength dependence imparted by dust (e.g., Whittet et al. 1992).
#### 2.3.3 Previous SN Ia Polarimetry Studies
Evidence for intrinsic polarization in SNe Ia was initially difficult to find, as the first investigations detected only marginally significant polarization among normal-brightness events observed near maximum, $`p<0.2`$% (e.g., Wang et al. 1996, 1997). Broad-band polarimetry of one overluminous event observed over a month after maximum also found no intrinsic polarization down to a level of $`p0.3\%`$ (Wang et al. 1996). Significant advances have been made in the last few years. Leonard et al. (2000b) reported the first convincing, albeit weak, features in the polarization of an SN Ia, SN 1997dt, which was likely a subluminous event; these data are presented and analyzed in more detail in the present study.
The first thorough spectropolarimetric study of an SN Ia is that by Howell et al. (2001). In spectropolarimetry obtained at maximum light, the subluminous SN 1999by exhibits a polarization change of $``$0.8% from 4800 Å to 7100 Å, and a sharp polarization modulation of $``$0.4% across the strong Si II $`\lambda 6355`$ absorption. These features are explained within the context of an ellipsoidal model with a global asphericity of $``$20%, observed equator-on. This physical picture was achieved by assuming the ISP to be the observed polarization at the blue edge of the spectrum, which was about $`0.2\%\mathrm{at}4800`$ Å. With this choice of ISP, the inferred intrinsic polarization rises from $`0\%`$ to about $`0.8\%`$ from blue to red, with a sharp depolarization from $`0.4\%`$ to near $`0\%`$ across the Si II $`\lambda 6355`$ feature. An argument supporting both the ellipsoidal model as well as the ISP choice is that after ISP removal the PA becomes nearly independent of wavelength. Howell et al. (2001) note, on the other hand, that a wavelength-independent PA also results if one assumes the ISP to be the observed polarization of the far red edge of the spectrum; in this case, the polarization modulation across the Si II $`\lambda 6355`$ becomes a polarization increase. Such a scenario could result from selective blocking of forward-scattered light, as described in § 2.3.1. However, Howell et al. (2001) find the theoretical arguments supporting the former ISP to be more compelling. With this ISP, the redward rise in intrinsic polarization is attributed to the decreasing importance of line opacities, and the increased influence of continuum electron scattering at longer wavelengths. The depolarization across the Si II $`\lambda 6355`$ absorption is also attributed to the depolarizing effect of line scattering. Despite the detection of intrinsic polarization for SN 1997dt and SN 1999by, their unusual (subluminous) nature nevertheless left some doubt about intrinsic polarization in normal SNe Ia.
That doubt has recently been put to rest with the work of Kasen et al. (2003) on SN 2001el. For this normal-luminosity event, the percent polarization changed from blue to red by $``$0.4% in spectropolarimetry obtained 1 week before maximum brightness. However, the extraordinary feature here is the existence of distinct high-velocity Ca II near-infrared (IR) triplet absorption ($`v`$ = 18,000–25,000 km s<sup>-1</sup>) in addition to the usual, lower-velocity Ca II feature. A similar, but much weaker, high-velocity feature had been previously observed in SN 1994D, and perhaps in other SNe Ia as well; the number of pre-maximum spectra covering the near-IR spectral range is small. The polarization is seen to change dramatically in this feature, by $``$0.4%. The Ca II near-IR feature is examined by Kasen et al. (2003) in an elegant study, which concludes that it is likely due to photospheric obscuration by a clumped shell of high-velocity material. Using multi-epoch data, Wang et al. (2003) demonstrate that the nature of the polarization changes over the course of two weeks following this early epoch, becoming nearly undetectable a week after maximum brightness, further solidifying the case for intrinsic polarization at early times. Although not commented on by either study, it appears that the Si II $`\lambda 6355`$ line also shows a polarization modulation in the earliest epoch, and with a different PA from the Ca II near-IR feature. Although it does not affect the main results, a cautionary note on the difficulty of ISP determination is set by the fact that the two studies arrive at quite different values: Kasen et al. (2003) derive an ISP of very nearly $`0\%`$ by assuming the blue edge of the spectrum to be unpolarized, whereas Wang et al. (2003) obtain an ISP of $``$0.6% by attributing the observed polarization $`38`$ days after maximum light, when SN 2001el is argued to be in the nebular phase, entirely to ISP.
Most recently, Wang et al. (2005a) analyze a single epoch of pre-maximum spectropolarimetry of SN 2004dt, an HV SN Ia for which an ISP of $`q_{\mathrm{ISP}}=0.2\pm 0.1\%`$ and $`u_{\mathrm{ISP}}=0.2\pm 0.1\%`$ is derived from the observed polarization of a handful of narrow, blue spectral regions in which no single spectral feature dominates in the total-flux spectrum. This results in rather low continuum polarization, $`p0.4\%`$, but very strong modulations across spectral lines. The polarization spikes reach $`2\%`$ in the deep troughs due to Si II $`\lambda 4130`$ and $`\lambda 6355`$; lesser peaks are observed in features identified with Mg II $`\lambda 4471`$, a blend of Si II $`\lambda \lambda 5041,5056`$ and Fe II $`\lambda \lambda 4913,5018,5169`$, and the Ca II near-IR triplet. All line polarization has similar directional behavior in the $`q`$$`u`$ plane, suggesting a common origin. Interestingly, whereas other strong line features in the total-flux spectrum are characterized by strong polarization modulations, O I $`\lambda 7774`$ shows no polarization signature. These features are explained in terms of optically thick bubbles of IMEs, the result of partial burning, that are asymmetrically distributed within an essentially spherical oxygen substrate that remains from the progenitor. Note that Wang & Wheeler (1997) find polarization variation at a level of $`>0.5\%`$ across the strong Si II $`\lambda 6355`$ feature of another HV SN Ia, SN 1997bp, in unpublished spectropolarimetry obtained near maximum light.
From this small sample, the observations thus far suggest that normal and, perhaps, overluminous events are weakly polarized ($`p0.4\%`$), with subluminous ones possessing somewhat larger values ($`p0.8\%`$). HV SNe Ia have modest continuum polarization, but possess the highest line polarization of all, achieving polarizations of up to $`2\%`$ in the strongest features. We assess the robustness of these tentative trends with the present study of four SNe Ia.
### 2.4 The Type Ia Supernovae 1997dt, 2002bf, 2003du, and 2004dt
SN 1997dt was discovered (Qiao et al. 1997) by the Beijing Astronomical Observatory Supernova Survey (Li et al. 1996) on 1997 November 22.44 (UT dates are used throughout this paper) at an unfiltered magnitude of $``$15.3 in the Sbc galaxy NGC 7448. Images of the same field taken eight days earlier show no star at the position of the SN to a limiting unfiltered magnitude of about 18.5. An optical spectrum obtained immediately (0.06 days) after discovery showed it to be a Type Ia event (Qiao et al. 1997); a subsequent examination by Li et al. (2001b) estimates the age at discovery to be $`7\pm 5`$ days relative to maximum light, which places the date of maximum at very roughly 1997 November 29. Tonry et al. (2003) report that SN 1997dt suffers from a host extinction of $`A_V=0.46\mathrm{mag}`$, the median of the values derived from MLCS and “Bayesian Adapted Template Match Method” (J. L. Tonry et al. 2005, in preparation; Tonry et al. 2003) analyses of its unpublished light curves. An MLCS2k2 fit (S. Jha et al., in preparation) suggests that it is subluminous, with $`\mathrm{\Delta }=0.94`$ (corresponding to a $`\mathrm{\Delta }m_{15}(B)1.8`$ mag; see § 2.1), although there is a long shoulder of probability to lower $`\mathrm{\Delta }`$ values, indicating that there is a wide range of light-curve shapes that can fit the sparse number of photometric points (2 in $`B`$, 3 in $`V`$, and 2 in $`I`$). From a pre-maximum spectrum posted at the Center for Astrophysics’ Recent Supernovae Page,<sup>2</sup><sup>2</sup>2http://cfa-www.harvard.edu/cfa/oir/Research/supernova/RecentSN.html we estimate $``$(Si II) $`0.3`$ (see Nugent et al. 1995), which implies $`\mathrm{\Delta }m_{15}(B)1.33`$ mag from the correlation derived by Benetti et al. (2004), also consistent with a somewhat subluminous classification.
SN 2002bf was discovered (Martin & Li 2002) by the Lick Observatory and Tenagra Observatory Supernova Searches (LOTOSS; Schwartz et al. 2000) on 2002 February 22.30 at an unfiltered magnitude of $``$17, very close to the nucleus of the Sb galaxy PGC 029953. An image of the same field taken twenty days earlier showed nothing at the position of SN 2002bf to a limiting unfiltered magnitude of $``$19. Optical spectra obtained on 2002 March 6.21 by Matheson et al. (2002) and on 2002 March 7.41 by Filippenko et al. (2002) identified it as a Type Ia event. Both groups noted that the expansion velocity derived from the absorption minimum of the Si II $`\lambda 6355`$ Å line was significantly greater than normal for an SN Ia near maximum light, indicating that it may be an HV SN Ia. A “spectral feature age” (Riess et al. 1997) of $`0\pm 2`$ days was derived from the March 6 spectrum (Matheson et al. 2002). The MLCS2k2 analysis of the light curves presented in § 4.2.1 yields 2002 March $`4.37\pm 0.50`$ as the date of maximum $`B`$ light, consistent with the age on March 6 derived by Matheson et al. (2002).
SN 2003du was discovered (Schwartz & Holvorcem 2003) by LOTOSS on 2003 April 22.4 at an unfiltered magnitude of $``$15.9 in the SBd galaxy UGC 9391. An image of the same field taken fifteen days before discovery showed nothing at the position of SN 2003du to a limiting unfiltered magnitude of $``$19. An optical spectrum obtained shortly thereafter, on 2003 April 24.06, identified it as a Type Ia event roughly two weeks before maximum light (Kotak et al. 2003). The MLCS2k2 analysis of the light curves presented in § 4.2.1 yields 2003 May $`6.12\pm 0.50`$ as the date of maximum $`B`$ light, consistent with the earlier spectral age estimate. It also agrees with the epochs of maximum determined by the recent photometric studies of SN 2003du by Anupama et al. (2005) and Gerardy et al. (2004).
SN 2004dt was discovered (Moore & Li 2004) by the Lick Observatory Supernova Search (Filippenko et al. 2003) at an unfiltered magnitude of $``$16.1 in the SBa galaxy NGC 799 on 2004 August 11.48. An image of the same field taken ten days earlier showed nothing at the position of SN 2004dt to a limiting unfiltered magnitude of $``$18. Spectra taken within 2 days of discovery by Gal-Yam (2004), Patat et al. (2004), and Salvo et al. (2004) confirmed it to be an SN Ia before maximum light. Patat et al. (2004) noted that several absorption lines showed high expansion velocities, a result confirmed by the recent study by Wang et al. (2005a), suggesting that, like SN 2002bf, SN 2004dt is an HV SN Ia. A preliminary analysis of the light curves of SN 2004dt shows that maximum light occurred near 2004 August 20 (W. Li, personal communication). A series of Hubble Space Telescope (HST) UV spectral observations was obtained as part of program GO-10182 (P.I. Filippenko), and will be analyzed in a future paper.
## 3 Observations and Reductions
We obtained single-epoch spectropolarimetry of SN 1997dt, SN 2002bf, SN 2003du, and SN 2004dt on days 21, 3, 18, and 4 (respectively) after maximum light. We also obtained additional optical spectroscopy and $`BVRI`$ photometry of SN 2002bf and SN 2003du. For SN 2002bf, our photometry samples $`10`$ to $`57`$ days from the time of maximum light, with one additional flux spectrum taken on day 9 after maximum. For SN 2003du, our ground-based photometry covers $`5`$ to $`113`$ days from maximum, with one additional epoch on day 436 taken using the High Resolution Channel (HRC) of the Advanced Camera for Surveys (ACS) on board HST. Five additional spectral epochs sample its development from days $`24`$ to $`82`$ after maximum.
### 3.1 Photometry
#### 3.1.1 Ground-Based Photometry of SN 2002bf and SN 2003du
All ground-based photometric data were obtained using either the 0.76-m Katzman Automatic Imaging Telescope (KAIT; Filippenko et al. 2001; Li et al. 2003) or the Nickel 1 m reflector (Li et al. 2001a), both located at Lick Observatory. Figures 1 and 2 show KAIT images of PGC 029953 and UGC 9391, the host galaxies of SN 2002bf and SN 2003du, respectively. Also labeled in the KAIT images are the “local standards” in both fields that were used to measure the relative SN brightness on non-photometric nights. We obtained 12 epochs of Johnson-Cousins $`BVRI`$ photometry (Johnson et al. 1966, for $`BV`$; Cousins 1981, for $`RI`$) for SN 2002bf, all taken with KAIT, and 38 epochs of $`BVRI`$ photometry for SN 2003du, 33 of them taken with KAIT and five with the Nickel telescope. We also obtained three pre-maximum unfiltered observations of SN 2002bf with KAIT (approximating the $`R`$ band, see Li et al. 2003), and one additional epoch of $`RI`$ photometry with KAIT of SN 2003du.
For the photometry we employed the usual techniques of galaxy “template” subtraction (Li et al. 2000, and references therein), point-spread function fitting (Stetson 1991, and references therein), and using “local standards” to determine the $`BVRI`$ brightnesses of the SNe on non-photometric nights; in general, we closely followed the technique detailed by Leonard et al. (2002d). We note that the galaxy subtraction procedure for SN 2002bf was particularly challenging since it is only $`4\stackrel{}{\mathrm{.}}1`$ from its host galaxy’s center.
The absolute calibration of the SN 2002bf field was accomplished on the photometric nights of 2002 May 14 and 2004 March 17 with the Nickel telescope, and 2003 February 3 and 2004 March 18 with KAIT, by observing several fields of Landolt (1992) standards over a range of airmasses in addition to the SN 2002bf field. The absolute calibration of the SN 2003du field was similarly derived from data taken on the photometric nights of 2003 May 31, June 1, June 26, and August 27 with the Nickel telescope, and of 2003 May 22 and 2004 March 18 with KAIT. The color terms used to transform the filtered instrumental magnitudes to the standard Johnson-Cousins system are those of Foley et al. (2003). We list the measured $`BVRI`$ magnitudes and the $`1\sigma `$ uncertainties, taken as the quadrature sum of a typical photometric error and the $`1\sigma `$ scatter of the photometric measurements from all of the photometric nights, of the local standard stars in Tables 1 and 2.
After deriving the $`BVRI`$ magnitudes of the SNe based on a comparison with each of the local standards, we took the weighted mean of the individual estimates as the final standard magnitude of the SNe at each epoch in each filter. The results of our ground-based photometric observations are given in Tables 3 and 4 and shown in Figures 3 and 4. The reported uncertainties come from the quadratic sum of the photometric errors (reported by DAOPHOT) and the transformation errors. For SN 2002bf, the uncertainty produced by the difficult galaxy-subtraction process contributed the majority of the error on nights with low S/N.
#### 3.1.2 Hubble Space Telescope Photometry of SN 2003du
We obtained HST images during the course of two orbits of a $`29^{\prime \prime }\times 26^{\prime \prime }`$ field of view centered on SN 2003du on 2004 July 15, 434 days after $`B_{\mathrm{max}}`$, with the ACS/HRC detector through filters F435W, F555W, F625W, and F814W (hereafter referred to as $`B,V,R,\mathrm{and}I`$, respectively), as part of our Snapshot survey program (GO-10272; P.I. Li) to investigate the late-time photometric behavior and environment of nearby SNe. SN 2003du was detected in all images. Total exposure times in $`BVRI`$ were, respectively, $`1680\mathrm{s}`$ (data archive designation j8z441011/3011), $`960\mathrm{s}`$ (j8z442011/4011), $`720\mathrm{s}`$ (j8z441021/3021), and $`1440\mathrm{s}`$ (j8z442021/4021).
To derive the HST photometry, we followed as closely as possible the procedure detailed by Sirianni et al. (2005), including correction for the effects of SN light contaminating the background region, aperture corrections, the “red-halo” effect (for the $`I`$-band), and CTE degradation (Riess 2004). We translated the resulting instrumental magnitudes to the standard Johnson-Cousins $`BVRI`$ system by using the coefficients and color corrections tabulated by Sirianni et al. (2005). The final results of our HST photometry are included in Table 4 and displayed in Figure 4.
### 3.2 Spectropolarimetry and Spectroscopy
We obtained single epochs of spectropolarimetry for SN 2002bf and SN 2003du on 2002 March 7 and 2003 May 24, respectively, with the Low-Resolution Imaging Spectrometer (Oke et al. 1995) in polarimetry mode (LRISp; Cohen 1996)<sup>3</sup><sup>3</sup>3Instrument manual available at http://www2.keck.hawaii.edu/inst/lris/pol\_quickref.html. at the Cassegrain focus of the Keck-I 10-m telescope. We observed SN 1997dt on 1997 December 20 with LRISp using the Keck II 10-m telescope, and SN 2004dt on 2004 August 24 with the Kast double spectrograph (Miller & Stone 1993) with polarimeter at the Cassegrain focus of the Shane 3-m telescope at Lick Observatory. We reduced the polarimetry data according to the methods outlined by Miller et al. (1988) and detailed by Leonard et al. (2001) and Leonard & Filippenko (2001).
The polarization angle offset between the half-wave plate and the sky coordinate system was determined by observing the following polarized standard stars from the list of Schmidt et al. (1992) and setting the observed $`V`$-band polarization position angle (i.e., $`\theta _V`$, the debiased, flux-weighted average of the polarization angle over the wavelength range 5050–5950 Å; see Leonard et al. 2001) equal to the cataloged value: BD $`+64^{}106`$ (1997 December 20), BD $`+59^{}389`$ (2002 March 7), and HD 161056 (2003 May 24). On the night of 2004 August 24, we averaged the polarization angle offsets derived from observations of three polarized standards from the Schmidt et al. (1992) list, HD 204827, BD $`+59^{}389`$, and HD 19820; the individual offsets were internally consistent to within $`1^{}`$. To check for instrumental polarization, the following null standards taken from the lists of Turnshek et al. (1990), Mathewson & Ford (1970), Schmidt et al. (1992), and Berdyugin et al. (1995), were also observed: HD 94851 (1997 December 20), HD 57702 (2002 March 7), HD 109055 and BD $`+32^{}3739`$ (2003 May 24), and HD 212311 (2004 August 24). All stars were measured to be null to within $`0.1\%`$, which is also our estimate of the systematic uncertainty of a continuum polarization measurement made with either the Keck or Lick spectropolarimeters (e.g., Leonard et al. 2001).
Additional specifics of the observations of SN 2002bf taken on 2002 March 7, including an investigation of the potential impact of second-order light contamination and instrumental polarization (both shown to be minimal) in the setup used on this night, are given by Leonard et al. (2002a). Second-order light contamination is not a concern for our observation of SN 1997dt due to its limited spectral range. For SN 2003du, the use of a dichroic (D560) to split the beam near 5600 Å eliminates second-order light contamination on the red side. Our spectropolarimetric observation of SN 2004dt was taken in a setting that included the use of an order-blocking filter (GG455) to prevent contamination by second-order light at red wavelengths. Beyond $``$9000 Å, however, second-order contamination may exist, but for reasons similar to those discussed by Leonard et al. (2002a) we believe it to have minimal impact for this particular object.
To derive the total-flux spectra, we extracted all one-dimensional sky-subtracted spectra optimally (Horne 1986) in the usual manner. Each spectrum was then wavelength and flux calibrated, and was corrected for continuum atmospheric extinction and telluric absorption bands (Wade & Horne 1988; Bessell 1999; Matheson et al. 2000). With the exception of the spectropolarimetric observations of SN 1997dt and SN 2002bf, all spectra were taken near the parallactic angle (Filippenko 1982), so the spectral shape should be quite accurate. Table 5 lists the spectropolarimetric and spectral observations for all four SNe. Figures 58 show the observed spectropolarimetry data of the four objects, and Figures 9 and 10 show the complete series of spectra obtained for SN 2002bf and SN 2003du, respectively.
## 4 Analysis
### 4.1 Spectroscopy
Our spectrum of SN 1997dt, taken $``$21 days after maximum, shows typical features for an SN Ia at this phase (Fig. 5a). Similarly, our spectral sequence of SN 2003du (Fig. 10) closely follows the evolution of the normal-luminosity SN Ia 1994D (Patat et al. 1996; Filippenko 1997) in terms of the strengths and blueshifts of line features. This is consistent with the analyses of Anupama et al. (2005) and Gerardy et al. (2004), in which spectra taken before and shortly after maximum light were also examined. This convincingly establishes SN 2003du as a spectroscopically “typical” SN Ia (Branch et al. 1993).
As discussed in § 2.4, the spectra of both SN 2002bf and SN 2004dt are peculiar in one regard: the blueshifts of many of the spectral lines, most noticeably Si II $`\lambda 6355`$, occur at significantly higher velocity than is typical for an SN Ia at this phase, and indicate that these are both HV SNe Ia.
For comparison, we have measured the velocities of the Si II $`\lambda 6355`$ line in a number of other HV SNe Ia from our database, along with the spectroscopically normal SN 1994D, the subluminous SN 1991bg (e.g., Filippenko et al. 1992a), and the overluminous SN 1991T (e.g., Filippenko et al. 1992b). We present the results in Table 6 and Figure 11, from which it is clear that both SN 2002bf and SN 2004dt belong to the class of HV SNe Ia; indeed, SN 2002bf is the most extreme HV SN Ia yet observed for its epochs. The figure also suggests that line velocity is not strongly correlated with luminosity, although the expansion velocity of the subluminous SN 1991bg is somewhat lower than typical values.
Figure 12 presents a spectral comparison near maximum light of three HV SNe Ia (SN 2002bf, SN 2002bo, and SN 2004dt) with the spectroscopically normal SN 1994D. The excessive blueshift of the Si II $`\lambda 6355`$ trough is obvious for the three HV SN Ia compared with SN 1994D. While it must be cautioned that the spectra span a range of ages of about five days near maximum light, a time when significant spectral development occurs, the blueshift differences are much greater than can be explained by age differences alone. In addition, the Si II line is significantly stronger in the HV SNe Ia compared with SN 1994D, with equivalent widths of $`140`$ Å compared with $``$100 Å for SN 1994D. Conversely, the O I $`\lambda 7774`$ absorption is relatively weaker in the HV SN Ia events, with equivalent widths of $`100`$ Å in the HV SNe Ia compared with $`125`$ Å for SN 1994D. The fact that the Si II line is significantly stronger in the HV SNe Ia sample, and the O I line relatively weaker, is consistent with the scenario in which a greater fraction of C and O is burned to IMEs in HV SN Ia than in more typical events; it also follows naturally from the GCD model, since the obscuring pancake is formed from the products of oxygen burning.
While HV SNe Ia share many spectral characteristics, it is clear that they also exhibit spectral diversity. For instance, whereas the absorption trough of the Ca II near-IR triplet is significantly blueshifted for SN 2002bf and SN 2002bo relative to SN 1994D ($`16,300\mathrm{km}\mathrm{s}^1`$ and $`14,500\mathrm{km}\mathrm{s}^1`$ for SN 2002bf and SN 2002bo, respectively, compared with $``$10,900 km s<sup>-1</sup> for SN 1994D, where we have assumed $`\lambda _0=8579`$ for the Ca II near-IR triplet, a value derived using the prescription given by Leonard et al. 2002c), it has a more normal blueshift ($`v=10,900`$ km s<sup>-1</sup>) in SN 2004dt. The equivalent widths of the Ca II near-IR lines also show a suggestive trend, with both SN 2002bf and SN 2002bo having widths of $`>200`$ Å, while both SN 2004dt and SN 1994D have equivalent widths of $``$100 Å. This may indicate that the explosive nucleosynthesis did not proceed up to Ca as far out in the atmosphere (or in the bubble in the GCD scenario) of SN 2004dt as it did in the other two HV SNe.
Finally, it is clear from Figure 12 that not all lines share the extreme velocities seen in the Si II $`\lambda 6355`$ feature. For instance, as first noticed by Wang et al. (2005a) in a pre-maximum spectrum of SN 2004dt, the S II “W” feature in the spectra of SN 2002bo and SN 2004dt indicates velocities comparable to those in SN 1994D; for SN 2002bf they are somewhat larger, but still not near the velocity of the Si II $`\lambda 6355`$ line. Wang et al. (2005a) propose that this may indicate that sulfur is more confined to the lower-velocity, inner region.
### 4.2 Photometry
#### 4.2.1 Ground-Based Photometry
The results of the MLCS2k2 application to the photometry of SN 2002bf and SN 2003du are given in Table 7. For details of the MLCS2k2 procedure used, see S. Jha et al. (in preparation); an overview of the technique is provided by Riess et al. (2005). The MLCS2k2 analysis finds SN 2002bf to be of typical luminosity. SN 2003du is slightly overluminous, although its pre-maximum spectral evolution (Anupama et al. 2005) demonstrates that it is not a SN 1991T-like event, as the strength of the Si II $`\lambda 6355`$ feature is comparable to that seen in spectra of normal SNe Ia.
For SN 2002bf, $`B_{\mathrm{max}}`$ occurred on 2002 March $`4.37\pm 0.50`$, which is $`10`$ days after discovery and 5 days before our filtered observations commenced. For SN 2003du, $`B_{\mathrm{max}}`$ occurred on 2003 May $`6.12\pm 0.50`$, which is $`14`$ days after discovery, and the same day that our $`BVRI`$ observations began. Our derived date of maximum light for SN 2003du agrees with those found by Anupama et al. (2005) and Gerardy et al. (2004) using independent data sets.
#### 4.2.2 HST Photometry of SN 2003du
Late-time photometry (e.g., $`t>200\mathrm{d}`$) of SNe Ia exists for only a handful of objects (see, e.g., Milne et al. 2001, and references therein). From the small sample, there are two main features of note. First, SN Ia decline rates are typically much faster than the decay slope of $`{}_{}{}^{56}\mathrm{Co}`$ $`{}_{}{}^{56}\mathrm{Fe}`$ of $`0.98\mathrm{mag}(100\mathrm{d})^1`$ predicts. This decay mechanism is thought to be primarily responsible for powering the luminosity from the early nebular phase out to $``$1000 days for SNe of all types. Essentially, the $`{}_{}{}^{56}\mathrm{Co}`$ decays release most of their energy in the form of $`\gamma `$-rays, which, given enough optical depth, can become trapped in the ejecta and Compton scatter off free electrons. The energetic electrons generate optical photons primarily through ionization and excitation of atoms, and the ejecta are transparent to these photons. A small portion ($``$3.5%, see Arnett 1979) of the total $`{}_{}{}^{56}\mathrm{Co}`$ decay energy comes in the form of positrons, which may deposit their kinetic energy in the ejecta and then annihilate with electrons, producing two $`\gamma `$-ray photons of energy $`E_\gamma =m_ec^2`$. The steeper decline seen in the late-time photometry of SNe Ia has been explained by significant transparency of the ejecta to $`\gamma `$-ray photons (Milne et al. 1999) and positron escape (Milne et al. 2001).
A quantitative measure of the late-time decline is given by Cappellaro et al. (1997), who investigate how the decline from maximum to 300 days after peak $`V`$ brightness, denoted $`\mathrm{\Delta }m_{300}(V)`$, correlates with intrinsic SN brightness as derived from the $`\mathrm{\Delta }m_{15}(B)`$ parameter. Although limited by small sample size (5 objects), they find a convincing correlation, with $`\mathrm{\Delta }m_{300}(V)`$ going from 6.7 mag to 8.4 mag as the sample runs from overluminous (SN 1991T) to subluminous (SN 1991bg) events, with the normal-brightness SN 1994D characterized by $`\mathrm{\Delta }m_{300}(V)=7.3`$ mag. A second feature that has thus far been seen in only two SN Ia events (SN 1991T and SN 1998bu; see Schmidt et al. 1994 and Cappellaro et al. 2001, respectively) is a sudden flattening of the late-time optical light curves, which has been attributed to the contribution of a light echo from foreground dust clouds.
Our photometric data from HST, taken $`436`$ days after $`B_{\mathrm{max}}`$, allow us to investigate the late-time photometric behavior of SN 2003du. From the inset of Figure 4, it is clear that SN 2003du, like all SNe Ia analyzed before it, declines significantly faster than the $`{}_{}{}^{56}\mathrm{Co}`$ $`{}_{}{}^{56}\mathrm{Fe}`$ decay rate. In fact, we measure the average decay rate in $`V`$ (shown by Milne et al. 2001 to track the bolometric luminosity of an SN Ia quite accurately) of SN 2003du between our last two photometric epochs on days $`113`$ and $`436`$ to be $`\mathrm{\Delta }V=1.47\pm 0.02\mathrm{mag}(100\mathrm{d})^1`$. The slope also appears to have been rather constant throughout the period between our two last epochs, as indicated by the good agreement of the late-time data taken from Anupama et al. (2005) near day 300 and the decay slope determined from our data alone. Using the Anupama et al. (2005) data point and our estimate of $`V_{\mathrm{max}}`$ (Table 7), we derive $`\mathrm{\Delta }m_{300}(V)=6.74`$, which is most consistent with the values found by Cappellaro et al. (1997) for overluminous events. That SN 2003du may be somewhat overluminous was also suggested by its peak $`V`$ magnitude of $`19.67\pm 0.02`$ mag, which is 0.17 mag brighter than the fiducial template used in the MLCS2k2 procedure (Table 2). There is no evidence from our data of any contribution to the SN brightness from a light echo, although we note that the major indications of additional contributions to the apparent brightness of SNe 1991T and 1998bu did not become obvious until epochs $`500\mathrm{d}`$.
### 4.3 Reddening
#### 4.3.1 Techniques to Estimate SN Ia Reddening
Accurate determination of SN reddening is crucial both for deriving intrinsic SN properties as well as interpreting spectropolarimetry, since the same dust that reddens SN light can also polarize it as discussed in § 2.3.2. When multi-band photometry is available, the MLCS2k2 technique can accurately estimate the total extinction (see Riess et al. 2005 for discussion). When this is lacking, other methods must be used.
Galactic extinction along the l-o-s is accurately estimated by the dust maps of Schlegel et al. (1998, hereafter SFD) to an estimated precision of $``$15%. For host-galaxy extinction, a rather crude approach, to which many investigators resort in the absence of photometry, is to employ the rough correlation found between the total equivalent width ($`W_\lambda `$) of the interstellar (IS) Na I D doublet ($`\lambda \lambda 5890,5896`$) and reddening (Barbon et al. 1990). The Barbon et al. correlation has been subsequently improved upon by Munari & Zwitter (1997), who derive a more precise relation that uses just the equivalent width of the Na I D2 ($`\lambda 5890`$) line. Both relations, however, warrant healthy degrees of skepticism since sodium is known to be only a fair tracer of the hydrogen gas column (especially in dense environments, where sodium may be heavily depleted; e.g., Cohen 1973), from which the dust column is then estimated. The dust-to-gas ratio also varies significantly among galaxies (e.g., Issa et al. 1990), and it could well be that in the case of SNe, circumstellar, rather than interstellar, dust is present, for which the dust-to-gas ratio (or the extinguishing properties of the dust itself) could be unusual. A final complication is that, at the resolution typical of most optical spectra, individual absorption components along the l-o-s, whose contributions should be considered separately to determine the total reddening, are blended into a single profile. In such situations, the Munari & Zwitter (1997) relation formally yields an upper limit to the reddening, and caution must be used since the derived value could seriously overestimate the actual reddening. Nonetheless, the Munari & Zwitter (1997) relation is still frequently used to get approximate reddening values, or upper limits, especially in the case of null detections of IS Na I D lines.
#### 4.3.2 The Reddening of the Four SNe Ia
For SN 2002bf and SN 2003du we shall adopt the MLCS2k2 values of $`E(BV)=0.08\pm 0.04`$ mag and $`E(BV)=0.01\pm 0.01`$ mag (respectively) reported in Table 7 (for $`R_V=3.1`$). For SN 1997dt, Tonry et al. (2003) report a host-galaxy reddening value of $`E(BV)_{\mathrm{Host}}=0.15`$ mag (assuming $`R_V=3.1`$), the median of the values derived from MLCS2k2 and “Bayesian Adapted Template Match Method” analyses. The SFD dust maps predict $`E(BV)_{\mathrm{SFD}}^{\mathrm{MW}}=0.06`$ mag for SN 1997dt, for a total estimate of $`E(BV)=0.21`$ mag.
Our spectrum of SN 1997dt also exhibits strong Na I D IS absorption at both the redshift of the host galaxy and the MW. The resolution of our spectrum, $``$5 Å (Table 5), while sufficient to deblend the Na I D doublet, is not fine enough to resolve the individual absorption components that likely contribute to the D1 and D2 profiles. For host-galaxy absorption, we measure $`W_\lambda ^{\mathrm{tot}}=0.77`$ Å, with $`W_\lambda ^{\mathrm{D2}}=0.42`$ Å and $`W_\lambda ^{\mathrm{D1}}=0.35`$ Å. This translates to $`E(BV)_{\mathrm{NaID}}^{\mathrm{Host}}0.21`$ mag from the Munari & Zwitter (1997) relation. For the MW, we measure $`W_\lambda ^{\mathrm{tot}}=0.76`$ Å, with $`W_\lambda ^{\mathrm{D2}}=0.44`$ Å and $`W_\lambda ^{\mathrm{D1}}=0.32`$ Å, yielding $`E(BV)_{\mathrm{NaID}}^{\mathrm{MW}}0.23`$ mag from the Munari & Zwitter relation. The values given by SFD and Tonry et al. (2003) are consistent with the upper limits derived from the sodium relation. The purported accuracy of the SFD Galactic reddening value and the upper limit set by the Munari & Zwitter relation for the host reddening would indicate an upper reddening limit of $`E(BV)_{\mathrm{total}}0.27`$ mag. This limit will prove to have important implications when we examine the spectropolarimetry of SN 1997dt in § 4.4.2. The fact that the Munari & Zwitter relation overpredicts the reddening for both the host galaxy and, especially, the MW, may indicate that multiple, unresolved components make up the IS line profiles. Given the upper reddening limit, we conclude $`E(BV)=0.21\pm 0.06`$ mag for SN 1997dt.
For SN 2004dt, $`E(BV)_{\mathrm{SFD}}^{\mathrm{MW}}=0.03`$ mag, and Wang et al. (2005a) report a range of $`E(BV)_{\mathrm{total}}=0.14`$ to 0.2 mag from analysis of the color of their pre-maximum spectrum. However, they also note that SN 2004dt does not show noticeable Na I D IS lines at the redshift of NGC 799. We confirm the lack of Na I D lines in our spectrum as well, although its poor resolution ($``$18 Å, see Table 5) makes deriving even an upper limit to the strength of the Na I D IS lines difficult. Formally, our procedure yields $`W_\lambda (3\sigma )=0.1`$ Å for host-galaxy Na I D, which translates to a predicted upper limit of $`E(BV)_{\mathrm{NaID}}^{\mathrm{Host}}<0.02`$ mag. We thus confirm the discrepancy noted by Wang et al. (2005a) between the lack of IS sodium absorption and reddening inferred from other methods. It could be that the spectroscopic peculiarities of HV SNe Ia produce photometric irregularities that affect the general reddening relations. As mentioned in § 2.2, Benetti et al. (2004) do find photometric peculiarities, albeit minor, in their study of another HV SN Ia, SN 2002bo. However, since the Na I D relation, especially for the host galaxy, is also prone to error, we have no convincing way to decide between the two estimates. We thus take the simple average of the low and high reddening estimates, and incorporate the disparity into our final estimate’s uncertainty. This yields $`E(BV)=0.11\pm 0.06`$ mag as our best reddening estimate for SN 2004dt.
### 4.4 Spectropolarimetry of Four SNe Ia
To summarize: SN 1997dt is likely a somewhat subluminous event that is thought to be reddened by $`E(BV)=0.21\pm 0.06`$ mag, with an upper limit of $`0.27`$ mag. SN 2002bf and SN 2004dt are HV SNe Ia, with SN 2002bf being the most extreme example yet observed. Finally, SN 2003du is a slightly overluminous event that is minimally reddened. At maximum light, a previous subluminous event has been found to be moderately polarized ($`p0.8\%`$), whereas normal to overluminous examples have been less so ($`p0.4\%`$). The two HV SNe Ia that have reported polarization measurements show the greatest polarization features yet observed for any SN Ia type, with values approaching $`2\%`$ in the strongest lines of SN 2004dt (§ 2.3.3).
We now examine our single-epoch spectropolarimetry of SN 1997dt, SN 2002bf, SN 2003du, and SN 2004dt obtained on days 21, 3, 18, and 4, respectively, after maximum light. For each SN, we shall first discuss the observed polarization and then attempt to remove the ISP through the technique of Wang et al. (2005a), which assumes that spectral regions lacking strong individual flux or polarization features at blue wavelengths ($`\lambda <5000`$ Å) are intrinsically unpolarized. Since the choice of the specific spectral regions to use for ISP determination is admittedly somewhat subjective, we shall be careful to point out how our conclusions would change with other ISP choices.
Following ISP removal and examination of the data in the $`q`$$`u`$ plane, we calculate the intrinsic polarization of the SN by rotating the $`q`$$`u`$ axes through a single angle, $`\theta `$, that places the greatest degree of polarization change across the spectrum along the rotated $`q`$ axis. The angle by which the axes are rotated is determined through a uniform-weight, least-squares fit to the ISP-subtracted data in the $`q`$$`u`$ plane. The polarization degree measured along the rotated $`q`$ axis is then referred to as the rotated Stokes parameter (RSP; Tran 1995; Leonard et al. 2001), whereas that measured along the rotated $`u`$ axis is denoted URSP.<sup>4</sup><sup>4</sup>4Note that RSP and URSP, as defined here, are identical to the Wang et al. (2001) definitions of the polarization along the “dominant” and “orthogonal” axes, $`P_d`$ and $`P_o`$, respectively. That is, by definition, RSP shows polarization strength variation along the chosen polarization angle (i.e., the dominant axis), whereas the URSP shows polarization strength variation along the axis that is orthogonal to the dominant axis in the $`q`$$`u`$ plane. The main difference between our construction of the RSP in the observed polarization plots (Figures 58) and the ISP-subtracted data (Figures 15, 17, and 19) is that in the former the polarization angle about which the RSP is determined is a smoothly varying function of wavelength (i.e., it is the polarization angle, $`\theta `$, smoothed over many bins; see Leonard et al. (2001) for more details on the procedure), whereas in the latter, RSP and URSP are determined with respect to a single, unchanging polarization angle. For consistency with spectropolarimetry work in other fields, as well as our own prior studies, we shall continue to use the RSP and URSP designations, rather than the $`P_d`$ and $`P_o`$ designations of Wang et al. (2001).
We shall then analyze the resulting spectropolarimetry within the context of the models described in § 2.3.1. By construction, the greatest degree of polarization change, especially from the continuum regions, should occur in the RSP. Examination of URSP is especially important in the line features, however, as it can be used to discriminate between the predictions of the ellipsoidal and clumpy-ejecta models. In particular, the simple ellipsoidal model demands the existence of a single value of $`\theta `$ capable of placing all line polarization changes along the RSP (i.e., URSP should show no polarization features).
#### 4.4.1 Two HV SNe Ia: SN 2002bf and SN 2004dt
We begin by considering the spectropolarimetry of the two HV SNe Ia, SN 2002bf and SN 2004dt, which are displayed in Figures 6 and 8, respectively. Both objects show very low levels of observed polarization, with $`p_V=0.03\%`$, $`\theta _V=62^{}`$ for SN 2002bf and $`p_V=0.25\%`$, $`\theta _V=146\mathrm{°}`$ for SN 2004dt, where $`p_V`$ and $`\theta _V`$ approximate the rest-frame $`V`$-band polarization and polarization angle derived by calculating the debiased, flux-weighted averages of $`q`$ and $`u`$ over the interval 5050–5950 Å (see Leonard et al. 2001). The polarization of SN 2004dt may exhibit an upward trend with wavelength, increasing from about $`0.3\%`$ at blue wavelengths to $``$1.0% at the red edge ($`\lambda =9600`$ Å); a similar but smaller trend exists in the data for SN 2002bf. Neither object shows significantly different polarization in the “continuum” region 6800–7800 Å from that observed in the heavily line-blanketed region below 5000 Å (§ 2.3.1), suggesting that the intrinsic continuum polarization is quite low for both objects.
The salient features of both data sets are the extraordinarily large polarization modulations ($``$2%) across certain P-Cygni lines, most notably the Ca II near-IR triplet for SN 2002bf and the Si II $`\lambda 6355`$ line for SN 2004dt. Smaller features may be discerned across other lines in both objects.
Before attempting ISP removal, it is instructive to compare our SN 2004dt data with those obtained by Wang et al. (2005a) taken eleven days earlier, when the SN was about seven days before maximum light. The interval separating the two observations is one marked by rapid spectral and photometric evolution. For SN 2004dt, the velocity of the minimum of the high-velocity Si II $`\lambda 6355`$ line recedes by about $`3700\mathrm{km}\mathrm{s}^1`$ (from $`17,200`$ to $`13,500\mathrm{km}\mathrm{s}^1`$), and the weaker, lower-velocity lines such as S II $`\lambda \lambda 5612,5654`$ ($`\lambda _0=5635`$ Å) decrease by about $`2100\mathrm{km}\mathrm{s}^1`$ (from $`11,000`$ to $`8900\mathrm{km}\mathrm{s}^1`$). Remarkably, in spite of the great changes that occur in the flux spectrum between the two epochs, we find the spectropolarimetry to be virtually unchanged. The overall polarization, in both magnitude and polarization angle, is consistent with the earlier epoch, and the detected line-polarization features, most notably those at Si II $`\lambda 6355`$ and the Ca II near-IR triplet, are also very similar. As with the Wang et al. data, our observations do not show significant polarization modulation across the important O I $`\lambda 7774`$ line, although this line is now much weaker in the total-flux spectrum than it was earlier. Weaker polarization features, which are clearly detected in the earlier epoch, are difficult to confirm in our lower S/N data, but our data are not inconsistent with the earlier measurements, within the errors.
As discussed in § 2.3.3, Wang et al. (2005a) explain the spectropolarimetry of SN 2004dt in terms of optically thick bubbles of IMEs that are asymmetrically distributed within an essentially spherical oxygen substrate that remains from the progenitor material. Wang et al. (2005a) further note that the polarization behavior of the high-velocity lines is similar to the low-velocity lines in the early epoch, implying that the structures that obscure the photosphere have great radial extent. Within this context, then, it may not be surprising that the polarization features have not evolved significantly between the two epochs, despite the great evolution of line velocity.
In order to directly compare the spectropolarimetry of SN 2002bf and SN 2004dt with each other, we must attempt to remove the ISP. To establish the ISP, we apply the technique of Wang et al. (2005a) and choose the spectral regions indicated in Figures 6a and 8a. For SN 2002bf, this yields $`(q_{\mathrm{ISP}},u_{\mathrm{ISP}})=(0.01\%,0.05\%)`$, or ISP$`{}_{\mathrm{max}}{}^{}=0.05\%`$ at $`\theta _{\mathrm{ISP}}=39\mathrm{°}`$ at an assumed peak wavelength of the ISP of $`\lambda _{\mathrm{max}}=5500`$ Å. For SN 2004dt we derive $`(q_{\mathrm{ISP}},u_{\mathrm{ISP}})=(0.3\%,0.2\%)`$, which is within the uncertainty of the ISP found by Wang et al. (2005a) through this same approach, $`(q_{\mathrm{ISP}},u_{\mathrm{ISP}})=(0.2\%,0.2\%)`$. In order to facilitate direct comparisons between our data and those of Wang et al. (2005a), we shall adopt their ISP value for our study of SN 2004dt as well.
After removal of the small amount of ISP for SN 2002bf, we derive a best-fitting axis with $`\mathrm{PA}=123\mathrm{°}`$ (Fig. 13), about which we calculate the intrinsic RSP and URSP shown in Figure 15. In a similar way, for SN 2004dt we derive a best-fitting axis of $`\mathrm{PA}=146\mathrm{°}`$, which may be compared with the PA of $``$150$`\mathrm{°}`$ found by Wang et al. (2005a). Again, since the two values do not differ significantly, we adopt the Wang et al. direction for ease of comparison between the two data sets. The axis and ISP choices for SN 2004dt are indicated in Figure 14, and the resulting RSP and URSP are shown in Figure 15.
From examination of Figure 15, it is clear that the two events have spectropolarimetric similarities. Both show low overall polarization and modulations across the Si II $`\lambda 6355`$ and Ca II near-IR triplet absorptions. The detailed character of the line polarizations do differ, however. For SN 2002bf there is a $`2\%`$ polarization change across the Ca II feature, with a more modest ($``$0.4%) feature detected in the Si II line. (When discussing overall “polarization change” across a line feature without specific reference to either the RSP or URSP, we effectively mean the quadrature sum of the changes seen in the two parameters.) For SN 2004dt, the situation is reversed. In the case of the Ca II line, this may be related to its relative strength and velocity in the two spectra, as both the equivalent width and velocity of the feature in the flux spectrum are much greater in SN 2002bf than in SN 2004dt. As discussed earlier (§ 4.1), it is possible that burning to Ca occurred more extensively in SN 2002bf than it did in SN 2004dt. The explanation for the Si II line disparity, however, is not so obvious, as the lines have similar strengths and velocities.
Clearly, strong line and weak continuum polarization levels disfavor the simple ellipsoidal model. Applying Eq. (1) to the data for SN 2002bf and SN 2004dt yields lower bounds on the expected continuum polarization of $`p_{\mathrm{cont}}0.7\%`$ and $`p_{\mathrm{cont}}1.3\%`$, respectively. For both objects, however, we measure $`p<0.4\%`$ in both the observed and ISP-subtracted data for the spectral region 6800–7800 Å, which is largely devoid of line opacity in SNe Ia near maximum (§ 2.3.1). The ellipsoidal model therefore appears to be ruled out as the explanation for the polarization characteristics of these objects.
The clumpy-ejecta model, on the other hand, provides a natural explanation for high line and low continuum polarization levels (§ 2.3.1). In particular, the GCD scenario can successfully explain (1) the high line velocities, (2) the large polarization change in the line troughs, (3) the great radial extent of the obscuring material (i.e., for $`m_{\mathrm{atm}}=0.08M_{}`$, the obscuring pancake spans the range 10,000–21,000 km s<sup>-1</sup>), and (4) the lack of significant continuum polarization. A remaining challenge is to explain the lack of any polarization change across the O I $`\lambda 7774`$ line, especially in the earlier Wang et al. (2005a) data for SN 2004dt, when this line is extremely strong in the total-flux spectrum. Recent nucleosynthesis calculations based on multi-dimensional (2D and 3D) hydrodynamical simulations of the thermonuclear burning phase in SNe Ia show that as much as $`40\%`$ to $`50\%`$ of the ejected matter in SNe Ia is unburned carbon and oxygen (Travaglio et al. 2004). This has led Wang et al. (2005a) to propose that the primordial oxygen is nearly spherically distributed, within which asymmetrically distributed IME clumps, or an absorbing pancake in the GCD model, are embedded.
With this in mind, an important prediction of the clumpy-ejecta model (including the GCD scenario) is that the O I $`\lambda 7774`$ line should, in fact, show a polarization change with a polarization angle that differs by $`90^{}`$ from what is observed in the Si II and Ca II lines since its distribution is essentially the inverse of these IMEs. However, some oxygen is probably also contained in the clumps or pancake, since it is produced by explosive carbon burning in small quantities and should be present when such burning products as magnesium exist, whose spectral signature is unequivocally seen in the spectra. This would tend to reduce the polarization level in the O I $`\lambda 7774`$ line. Neither our data, nor those of Wang et al. (2005a), are of sufficiently high S/N to detect such a change, and it must be left to future, higher-S/N studies focused especially at early times when the O I $`\lambda 7774`$ feature is strong. Finding such a PA change would further strengthen the case for the clumpy-ejecta and/or GCD scenario.
We note that the ejecta-hole model of Kasen et al. (2004) is also capable of producing large line polarization with weak continuum polarization for sight-lines near to the hole (§ 2.3.1). Arguing against this model in these cases, though, is that it offers no natural explanation for why high line velocities should be associated with sight lines near to the hole. In fact, the Kasen et al. models predict lower absorption velocities when viewing down the hole.
In all, then, our study finds that HV SNe Ia are, as a group, characterized by much stronger line-polarization features than are seen in other SN Ia varieties. The ellipsoidal model is incapable of explaining the polarization characteristics of these objects, whereas the clumpy-ejecta and ejecta-hole models are more successful. On balance, the case for clumpy ejecta appears to be the most convincing explanation, with the GCD model investigated by Kasen & Plewa (2005) able to reproduce many of the observed spectral and spectropolarimetric features.
#### 4.4.2 SN 1997dt
We next turn to the likely subluminous SN 1997dt, which was observed about three weeks past maximum light. Figure 5 reveals an extraordinarily high level of observed polarization, $`p_V=3.46\%`$ at $`\theta _V=112^{}`$. This is by far the largest polarization yet observed for an SN Ia. A distinct polarization feature is detected in the Fe II $`\lambda 4555`$ trough and probably also in the Si II $`\lambda 5972`$ \+ Na I D and Si II $`\lambda 6355`$ lines. The degree of change in the Fe II $`\lambda 4555`$ feature reaches nearly $`1\%`$ in the $`q`$ parameter.
This amount of observed continuum polarization is surprising, given that our best total reddening estimate predicts an upper bound on the ISP of only $`1.89\%`$ from the Serkowski et al. (1975) relation (see § 2.3.2). If we trust both the reddening estimate and the upper ISP bound, then an intrinsic SN polarization of at least $`1.57\%`$ must exist to explain the observed polarization, far higher than has been indicated for any previous SN Ia. However, the technique of Wang et al. (2005a) suggests a much higher ISP level: $`(q_{\mathrm{ISP}},u_{\mathrm{ISP}})=(2.53\%,2.68\%)`$, or ISP$`{}_{\mathrm{max}}{}^{}=3.60\%`$, $`\theta _{\mathrm{ISP}}=113\mathrm{°}`$, for $`\lambda _{\mathrm{max}}=6500`$ Å, the wavelength of maximum ISP that yields the most convincing Serkowski-law fit to the observed polarization (Fig. 5d). If this truly is the ISP, then it implies an extraordinarily high polarization efficiency for the dust along the l-o-s to SN 1997dt: $`\mathrm{ISP}/E_{BV}18\%\mathrm{mag}^1`$, which is double the observed Galactic limit of $`9\%\mathrm{mag}^1`$. Of course, some of this discrepancy could be removed if the true reddening were greater than we suspect. However, even allowing the reddening to equal the upper limit of $`E(BV)_{\mathrm{total}}<0.27`$ mag set in § 4.3.2 still requires the dust-polarization efficiency to exceed the Galactic limit.
It thus appears that we face a stark choice: either SN 1997dt has the highest intrinsic polarization of any SN Ia yet observed, or the dust along the l-o-s has an exceptionally high polarization efficiency. The epoch of our observation of SN 1997dt is unique for a subluminous SN Ia, so there is no empirical database from which to draw expectations and help decide between these two options.
The situation for SN 1997dt, while perplexing, in fact is not unique: a similar palette of possibilities presented themselves in our previous study of a single epoch of spectropolarimetry of SN 1999gi, an SN II-P that also had a low reddening and large observed polarization (Leonard & Filippenko 2001; Leonard et al. 2002d). Like SNe Ia, SNe II-P as a group have historically shown very low intrinsic continuum polarization. After considering many polarization production mechanisms, including polarization due to newly formed dust in the SN ejecta and dust reflection by one or more off-center dust blobs external to the SN, we concluded that the most likely explanation for the polarization of SN 1999gi was that the host-galaxy dust along the l-o-s possesses a very high polarization efficiency, $`\mathrm{ISP}/E[BV]=31_9^{+22}\%\mathrm{mag}^1`$, which remains the largest value yet inferred for a single sight line in either the MW or an external galaxy (Leonard et al. 2002b).
There are arguments favoring a similar explanation here. First, a Serkowski law reasonably fits the observed continuum polarization of SN 1997dt (Fig. 5d). Further, at 21 days past maximum, SN 1997dt may be nearing the end of its photospheric phase, a time when spectropolarimetry may be losing its efficacy as an asymmetry probe due to a lack of electrons available to scatter the light. If we believe this to be the case, then we should not expect large polarization, even if the SN is highly aspherical. If we accept the large ISP, then we naturally would like to know whether it is due to dust in the MW or NGC 7448. Unfortunately, we have not observed any distant Galactic “probe” stars (§ 2.3.2) near to the l-o-s to estimate the Galactic ISP. There are, however, reasons to suspect that it is low. First, 20 stars within $`10\mathrm{°}`$ of the l-o-s are listed in the agglomeration of stellar polarization catalogs by Heiles (2000), and the greatest observed polarization is only $`0.3\%`$. Second, the great majority of the reddening is due to host-galaxy dust, since $`E(BV)_{\mathrm{Host}}=0.15`$ mag while $`E(BV)_{\mathrm{MW}}=0.06`$ mag (§ 4.3.2). As was the case with SN 1999gi, we would again conclude that it is the dust within the host galaxy that must have the extraordinarily high polarization efficiency.
There are arguments to oppose this, however. Previous SN Ia polarization studies have found intrinsic polarizations increasing toward red wavelengths (§ 2.3.3), which could certainly mimic a Serkowski law over the limited wavelength band covered by our observations. Further, the line polarization features demonstrate that the SN must possess at least some intrinsic polarization. Finally, the relatively late phase of the observation, invoked previously to argue against high intrinsic polarization, can also be used to argue in favor of it: at the stage immediately before an SN begins the transition to the nebular phase, the deepest layers of the ejecta are revealed. Although in need of confirmation by detailed modeling, for this epoch one can plausibly argue that an optical photosphere still exists with sufficient optical depth to electron scattering, perhaps even reaching the single-scattering limit, which is the most polarizing atmosphere (Höflich 1991). If the explosion mechanism itself is asymmetric, the largest imprint of the asymmetry would presumably be in the innermost layers, which could lead to a very large intrinsic polarization that reveals itself just at this late phase.
In fact, such an effect is seen in the spectropolarimetry of the SN II-P 2004dj as it transitions to the nebular phase (D. C. Leonard et al., in preparation). However, in the case of SN 2004dj, the strong increase is observed primarily in the “continuum” region 6800–7800 Å (and in a few strong line troughs), not in the overall level across the whole spectrum. Although in need of confirmation by detailed modeling, depolarizing line blanketing may also be significant for SNe Ia at blue wavelengths at this epoch, which would again argue for a large ISP as the explanation of the high observed polarization. One possibility that circumvents this difficulty is that asymmetrically distributed concentrations of radioactive Ni, recently uncovered in the thinning ejecta at this relatively late epoch, are responsible for the large polarization.
Curiously, the single, high-polarization observation of SN 1999gi occurred at a similar stage of its evolution, right at the end of the optical plateau that characterizes the photospheric phase in SNe II-P. It is thus unfortunate that no other spectropolarimetric epochs were obtained for either event (SN 1999gi and SN 1997dt) to serve as a basis for comparison. Multi-epoch data covering the transition from the photospheric to the nebular phases of SNe of all types will certainly help reveal more definitively the physical explanation for such high observed polarizations at these late times.
It may be possible to gain insight into the cause of the observed polarization from examination of the spectropolarimetry data in the $`q`$$`u`$ plane, shown in Figure 16. ISP originating from a single source (i.e., characterized by a single PA) will spread intrinsically unpolarized data points along a line in a direction that intersects the origin in the $`q`$$`u`$ plane. Finding that data lie predominantly along such a line, and exhibit a Serkowski law spectral shape, provides compelling evidence that a large, single source of ISP dominates the observed signal. Such was the case for SN II-P 1999gi (Leonard et al. 2002b). Similarly, for SN 1997dt, an elongation of the data points from blue to red wavelengths in a direction that roughly points back toward the origin also exists.
When the ISP derived earlier (ISP$`{}_{\mathrm{max}}{}^{}=3.60\%`$, $`\theta _{\mathrm{ISP}}=113\mathrm{°}`$, for $`\lambda _{\mathrm{max}}=6500`$ Å) is removed, the wavelength dependence of the polarization largely disappears (Fig. 16), strengthening the argument for a single, dominant ISP source. The data are then seen to lie along an axis with a fairly well-defined PA of $`58\mathrm{°}`$ (Fig. 16), against which we calculate the RSP and URSP shown in Figure 17. Although of individually low to moderate significance, the polarization features in the Fe II $`\lambda 4555`$ and Si II $`\lambda 6355`$ troughs, seen in both RSP and URSP, do seem to prefer the same general direction, with most of the modulation occurring along the URSP axis (e.g., the direction perpendicular to the main axis in the $`q`$$`u`$ plane). Taken at face value, consistent line-trough polarization changes in a direction differing from that preferred by the continuum favors an origin in the selective blocking of photospheric light by clumpy and asymmetrically distributed intermediate and iron-peak elements overlying the photosphere as opposed to an ellipsoidal scenario.
Other plausible ISP choices, however, yield different conclusions. For instance, a smaller ISP of $`2.6\%`$ at the same PA yields a more nearly spherical constellation of points centered near $`(q,u)=(0.7\%,0.07\%)`$, with the line excursions now pointing back toward the origin. Since some modelers predict that, in SNe Ia with an ellipsoidal asphericity, polarization decreases may exist in absorption troughs (§ 2.3.1), the ellipsoidal model cannot therefore be ruled out in this case as the cause of the polarization. Given the marginal significance of the weaker features, additional interpretation of the line-trough polarization degrees and directions is probably not warranted with these data.
In conclusion, we find evidence for a large ISP contribution to the observed polarization of SN 1997dt, probably in the range $`2.6\%\mathrm{ISP}3.6\%`$. This implies a polarization efficiency for the dust along the l-o-s in NGC 7448 that exceeds the Galactic limit. We do, however, also find evidence for polarization intrinsic to the object, most convincingly in specific line features and, perhaps, in the continuum as well. A combination of asymmetrically distributed radioactive Ni and synthesized IMEs overlying the photosphere may provide the simplest explanation for the line and potential continuum polarization at this late phase, although the ellipsoidal model cannot be definitively ruled out.
#### 4.4.3 SN 2003du
Our spectropolarimetry of SN 2003du represents the highest S/N data of our study, and permits a more detailed analysis of line features than has been possible for the other objects. SN 2003du presents a low level of observed polarization across most of the spectrum, with $`p_V=0.04\%`$, $`\theta _V=17^{}`$ (Fig. 7). The polarization exhibits an increasing trend with wavelength, rising from nearly zero at blue wavelengths to $``$0.2% at the red edge of the spectrum. There are distinct and significant polarization changes across several absorption features, including the Ca II near-IR triplet, the Si II $`\lambda 6355`$ line, probably a few weaker lines such as Fe II $`\lambda 4924`$, and very tentatively the Ca II H & K absorption.
Applying the technique of Wang et al. (2005a), and choosing the regions indicated in Figure 7a to estimate the ISP, yields $`(q_{\mathrm{ISP}},u_{\mathrm{ISP}})=(0.02\%,0.0\%)`$, or ISP$`{}_{\mathrm{max}}{}^{}=0.02\%`$, $`\theta _{\mathrm{ISP}}=90\mathrm{°}`$, for an assumed $`\lambda _{\mathrm{max}}=5500`$ Å. This very low ISP is consistent with the negligible reddening found earlier. Furthermore, on the same night the SN 2003du data were taken, we observed the distant Galactic star BD $`+50^{}1593`$ (spectral type F8, $`V=10.64`$ mag), located just $`0.28^{}`$ from the l-o-s of SN 2003du, and found it to be null to within $`0.1\%`$. From its spectroscopic parallax, we estimate BD $`+50^{}1593`$ to be at least $`190`$ pc away which, at a Galactic latitude of $`53^{}`$, satisfies the criterion of Tran (1995) that a good MW “probe” star be more than 150 pc from the Galactic plane. We thus have multiple reasons to suspect little ISP contaminating the data.
After removal of the minimal ISP, a well-defined axis with $`\mathrm{PA}=107\mathrm{°}`$ is derived (Fig. 18), about which we calculate the intrinsic RSP and URSP, shown in Figure 19. Compared with what was seen in the HV SNe Ia, and even SN 1997dt, the line-polarization features are not large, amounting to no more than $`0.3\%`$. However, the very high S/N of these data makes the detections unequivocal, and establishes intrinsic polarization in a spectroscopically and photometrically normal SN Ia at the latest phase yet observed.
The behavior of the line-trough polarization for the Ca II near-IR triplet and Si II $`\lambda 6355`$ are very similar. Both show sharp increases of $``$0.2% in RSP, as well as overall increases in URSP. Modest RSP increases of $`0.1\%`$ may also be discerned in the Fe II $`\lambda 4924`$ and Ca II H & K absorptions. At this phase, the O I $`\lambda 7774`$ feature in SNe Ia is quite weak and, coupled with telluric A-band contamination, makes definitive detection of polarization modulation in this important region difficult; the observed changes are at about the level of the statistical noise.
The commonality between the polarization behavior of the Si and Ca lines argues for similar origins. Given the high S/N of the data, we can examine rather closely the basic predictions of the ellipsoidal model that (a) the polarization angle should be independent of wavelength, and (b) the overall polarization should increase from blue to red wavelengths, with an expectation of $`p0\%`$ at $`\lambda 5000`$ Å (if the line opacity remains strong in these regions at the epoch of observation; § 2.3.1).
With our initial choice of ISP, the first criterion is clearly not satisfied, as there are obvious excursions in URSP throughout the spectrum, with especially sharp modulations occurring across the strongest lines. In fact, for ISP values constrained to lie along the symmetry axis we are unable to find any value that convincingly satisfies both criteria of the ellipsoidal model: values in the upper-right quadrant (e.g., $`[q_{\mathrm{ISP}},u_{\mathrm{ISP}}]=[0.3\%,0.18\%]`$) produce a polarization decrease with wavelength across the spectrum, while ISP values in the lower-left quadrant that satisfy the criterion of $`p0\%`$ at blue wavelengths (e.g., $`[q_{\mathrm{ISP}},u_{\mathrm{ISP}}]=[0.1\%,0.07\%]`$) result in strong PA changes across the spectral lines. Choosing an ISP point far beyond the constellation of data points, which may serve to make the PA changes less objectionable (although still statistically significant), has the unfortunate consequence of straining the limits implied from the very low reddening, at least for dust with normal polarizing efficiency.
In addition, examining the URSP behavior of the Si II $`\lambda 6355`$ and Ca II near-IR lines in Figure 19 more closely, we see that the generally increasing trends in URSP across the lines show sharp decreases right at the locations of peak RSP modulation, although the statistical significance of the modulations, especially for Ca II, is not high. If the abrupt changes in URSP are real, such structure is readily explained under the clumpy-ejecta model by variations in the distribution of the IMEs as a function of radius in the expanding ejecta. Chugai (1992) also demonstrates that such sharp changes in line features can be produced by excitation asymmetry. On the other hand, there is no obvious mechanism to produce such an effect in the simple ellipsoidal models. We thus conclude that it is difficult to reconcile the basic predictions of the ellipsoidal model with the data for SN 2003du.
We therefore suspect either clumps in the ejecta overlying the photosphere or ionization asymmetry as the cause of the inferred line and, perhaps, continuum polarization, and are led again to disfavor the ellipsoidal model as the cause of the intrinsic polarization of this SN.
## 5 Conclusions
We present post-maximum single-epoch spectropolarimetry of four SNe Ia, bringing to six the number of SNe Ia thus far examined in detail with spectropolarimetry during the early phases. The four objects span a range of spectral and photometric properties, yet all are demonstrated to be intrinsically polarized. This suggests that asphericity and/or asymmetry may be a ubiquitous characteristic of SNe Ia in the first weeks after maximum light. The nature and degree of the polarization varies considerably within the sample, but in a way that is consistent with, and extends, previously suspected trends. Our main spectropolarimetry results are as follows:
1. SN 2002bf and SN 2004dt, both HV SNe Ia observed shortly after maximum brightness, exhibit the largest polarization features yet seen definitively for any subtype of SN Ia, with modulations of up to $``$2% in the troughs of the strongest lines. The overall polarization level of both objects is minimal, at least at blue wavelengths; there is a possible trend of increasing polarization with wavelength for both objects, though neither shows significant polarization in the “continuum” region 6800–7800 Å. The ISP contamination is not thought to be large in either object.
2. SN 1997dt, believed to be a somewhat subluminous event, has the highest observed overall polarization of any SN Ia yet studied, $`p_V=3.46\%`$, at 21 days past maximum light. This demands either an extraordinarily large polarization efficiency for the dust along the l-o-s in NGC 7448, the largest intrinsic SN Ia polarization thus far found, or perhaps some combination of the two. The observed polarization rises by about $`0.5\%`$ from blue ($`\lambda =4300`$ Å) to red ($`\lambda =6700`$ Å) wavelengths, approximating a Serkowski-law ISP curve rather convincingly, albeit one with a somewhat unusual peak wavelength ($`\lambda _{\mathrm{max}}6500`$ Å). A polarization modulation of nearly $`1\%`$ in the strong Fe II $`\lambda 4555`$ absorption, and a more modest change of $`0.3\%`$ in the Si II $`\lambda 6355`$ line, demonstrate that the SN does possess intrinsic polarization features. However, we conclude that ISP is responsible for the bulk of the overall polarization that is observed, with $`2.6\%\mathrm{ISP}3.6\%`$, and that the polarization efficiency of the dust along the l-o-s in NGC 7448 likely exceeds the empirical Galactic limit.
3. SN 2003du is a slightly overluminous SN Ia. Our spectropolarimetry, taken 18 days after maximum light, is the highest S/N data obtained for our sample of objects. It reveals a low continuum polarization that increases by $``$0.3% from blue to red wavelengths, with distinct changes of $``$0.2% detected in the Si II $`\lambda 6355`$ and Ca II near-IR triplet lines; smaller changes may be detected in weaker lines. The very similar behavior of the polarization in the two strongest lines, in both magnitude and direction in the $`q`$$`u`$ plane, suggests a common polarization origin. ISP is thought to be minimal.
Ordered by increasing strength of line-polarization features in SNe Ia, we find as follows: ordinary/overluminous $`<`$ subluminous $`<`$ HV SNe Ia, with the strength of the line-polarization features increasing from $`0.2\%`$ in the slightly overluminous SN 2003du to $`2\%`$ in both HV SNe Ia in our study. Absolute continuum polarization levels are more difficult to establish, due largely to uncertainties in the ISP, but there are compelling reasons to believe that at least three of our objects possess very little intrinsic polarization in spectral regions outside of specific line features. The Howell et al. (2001) study of SN 1999by and our data on SN 1997dt provide some evidence that continuum polarization may be higher in subluminous objects than in other types.
There are a number of alternatives from which to choose for the origin of polarization in SNe Ia, including global asphericity (e.g., the ellipsoidal model), ionization asymmetry, and clumps in the ejecta overlying the photosphere. The small, redward rise in the overall polarization level that is discerned in at least three of our objects can be reproduced by models possessing either global asphericity or an ionization asymmetry. Under the ellipsoidal model, the levels of continuum polarization for SN 2002bf, SN 2003du, and SN 2004dt imply minor to major axis ratios of around $`0.9`$ if viewed equator-on (Höflich 1991; Wang et al. 2003). This level of asphericity would produce a luminosity dispersion of about 0.1 mag for random viewing orientations (Höflich 1991), which could explain some of the dispersion seen in the brightness-decline relation of SNe Ia. If the proposed global asphericity is more complicated, then the luminosity of a single SN Ia may depend on viewing angle in a non-trivial way such that, even for a large sample of objects, an overall bias to slightly higher or lower values may result.
The potential for significantly greater continuum polarization, perhaps of $``$1%, in the likely subluminous SN 1997dt observed $``$21 days after maximum would imply a more severe distortion, of at least $`20\%`$, from the models of Höflich (1991). Since most cosmological applications of SNe Ia rely on data acquired closer to maximum light, such a late-time asphericity, even if common, would not seriously affect the utility of SNe Ia as distance indicators.
To explain the ubiquitous line polarization, the simple ellipsoidal model is effectively ruled out for three of our objects, including, most convincingly, the two HV events. From a number of lines of reasoning, the most convincing explanation is partial obscuration of the photosphere by clumps of newly synthesized IMEs forged in the explosion.
For the HV SNe Ia, in particular, the GCD model studied by Kasen & Plewa (2005) provides a plausible explanation for many of the observed spectral and spectropolarimetric characteristics. It predicts the existence of an optically thick pancake of material with significant radial extent that partially obscures the optical photosphere, producing larger line velocities and equivalent widths for many spectral features, and stronger line-trough polarization than is seen in more typical events. These qualitative expectations are borne out by our data: The line features of SN 2002bf and SN 2004dt possess the strongest polarization modulations and greatest equivalent widths of our sample. The astonishing similarity between our epoch of spectropolarimetry of SN 2004dt, $``$4 days after maximum, and that presented by Wang et al. (2005a) from 11 days earlier, provides compelling evidence that the obscuring material also possesses great radial extent in the thinning ejecta.
That SNe Ia may be separable into different groups based on their spectropolarimetric characteristics yields one more clue to assist in narrowing down progenitor possibilities and/or models for the physics of the explosion. The assertion that some fraction of the IMEs in the ejecta of SNe Ia may be confined to bubbles or filaments is, however, a rather blunt discriminatory tool: At the present stage of theoretical modeling, deflagration, delayed-detonation, off-center delayed-detonation, and GCD models can all plausibly be argued to produce clumps in the ejecta (e.g., Wang et al. 2005a, and references therein). The specific predictions of the GCD model in particular need to be further examined, preferably in full three-dimensional simulations, to test whether the explosion mechanism itself is viable, and, if so, whether it can quantitatively reproduce the observed characteristics of at least some SNe Ia in detail. On the observational front, higher S/N data, preferably obtained at multiple epochs, will help to narrow down the possibilities as well. With the steady advances being made in the theoretical understanding of these events, and the growing rate of SNe Ia studied in detail with spectropolarimetry, prospects for improving our understanding of these events are bright.
We thank Aaron Barth, Louis-Benoit Desroches, Mohan Ganeshalingam, Deborah Hutchings, Ed Moran, and Karin Sandstrom for assistance with some of the observations and data reduction, and Saurabh Jha for producing MLCS2k2 fits for two of our objects. We thank an anonymous referee for useful suggestions that resulted in an improved manuscript. This research has made use of the NASA/IPAC Extragalactic Database (NED), which is operated by the Jet Propulsion Laboratory, California Institute of Technology, under contract with NASA. The work of A.V.F.’s group at UC Berkeley is supported by National Science Foundation (NSF) grant AST-0307894. D.C.L. is supported by an NSF Astronomy and Astrophysics Postdoctoral Fellowship under award AST-0401479. Additional funding was provided by NASA grants GO-9155, GO-10182, and GO-10272 from the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS 5-26555. A.V.F. is grateful for a Miller Research Professorship at UC Berkeley, during which part of this work was completed. Some of the data presented herein were obtained at the W. M. Keck Observatory, which is operated as a scientific partnership among the California Institute of Technology, the University of California, and NASA; the Observatory was made possible by the generous financial support of the W. M. Keck Foundation. KAIT was made possible by generous donations from Sun Microsystems, Inc., the Hewlett-Packard Company, AutoScope Corporation, Lick Observatory, the NSF, the University of California, and the Sylvia & Jim Katzman Foundation. The assistance of the staffs at Lick and Keck Observatories is greatly appreciated. |
warning/0506/math0506078.html | ar5iv | text | # Tannakian duality for Anderson-Drinfeld motives and algebraic independence of Carlitz logarithms
## 1. Introduction
### 1.1. Periods of $`t`$-motives
#### 1.1.1. Notation
Let $`𝔽_q`$ be the field of $`q`$ elements, where $`q`$ is a power of a prime $`p`$. Let $`k:=𝔽_q(\theta )`$, where $`\theta `$ is transcendental over $`𝔽_q`$, and define an absolute value $`||_{\mathrm{}}`$ at the infinite place of $`k`$ so that $`\left|\theta \right|_{\mathrm{}}=q`$. Let $`k_{\mathrm{}}:=𝔽_q((1/\theta ))`$ be the $`\mathrm{}`$-adic completion of $`k`$, let $`\overline{k_{\mathrm{}}}`$ be an algebraic closure, let $`𝕂`$ be the $`\mathrm{}`$-adic completion of $`\overline{k_{\mathrm{}}}`$, and let $`\overline{k}`$ be the algebraic closure of $`k`$ in $`𝕂`$.
#### 1.1.2. Anderson $`t`$-motives
Let $`t`$ be a variable over $`𝔽_q`$ that is independent from $`\theta `$, and let $`\overline{k}[t;𝝈]`$ be the ring of polynomials in $`t`$ and $`𝝈`$ over $`\overline{k}`$ subject to the relations
$$ct=tc,𝝈t=t𝝈,𝝈c=c^{1/q}𝝈,c\overline{k}.$$
An Anderson $`t`$-motive is a left $`\overline{k}[t;𝝈]`$-module $`𝖬`$ that is free and finitely generated as both a left $`\overline{k}[t]`$-module and as a left $`\overline{k}[𝝈]`$-module and that satisfies $`(t\theta )^n𝖬𝝈𝖬`$ for all $`n`$ sufficiently large (see §3.4). Anderson $`t`$-motives were originally defined in , where they were called “dual $`t`$-motives.”
#### 1.1.3. Rigid analytic triviality
We let $`𝕋:=𝕂\{t\}`$ be the Tate algebra of power series in $`𝕂[[t]]`$ that are convergent on the closed unit disk in $`𝕂`$, and let $`𝕃𝕂((t))`$ be its fraction field. Let $`𝔼`$ be the subring of $`𝕋`$ consisting of power series that are everywhere convergent and whose coefficients lie in a finite extension of $`k_{\mathrm{}}`$. Finally, for a Laurent series $`f=_ia_it^i𝕂((t))`$ and an integer $`n`$, we set $`\sigma ^n(f):=f^{(n)}:=_ia_i^{q^n}t^i`$.
If $`𝖬`$ is an Anderson $`t`$-motive and $`𝗆\mathrm{Mat}_{r\times 1}(𝖬)`$ has entries comprising a $`\overline{k}[t]`$-basis of $`𝖬`$, then there is a matrix $`\mathrm{\Phi }\mathrm{Mat}_r(\overline{k}[t])`$ representing multiplication by $`𝝈`$ on $`𝖬`$ so that
$$𝝈𝗆=\mathrm{\Phi }𝗆$$
and $`det\mathrm{\Phi }=c(t\theta )^s`$ for some $`c\overline{k}^\times `$ and $`s1`$. The Anderson $`t`$-motive is rigid analytically trivial (see Proposition 3.4.7) if there is a matrix $`\mathrm{\Psi }\mathrm{GL}_r(𝕋)`$ so that
$$\mathrm{\Psi }^{(1)}=\mathrm{\Phi }\mathrm{\Psi }.$$
It can be shown that the entries of $`\mathrm{\Psi }`$ are in fact in $`𝔼`$ (see Proposition 5.1.3).
#### 1.1.4. Connection with $`t`$-modules
The category of rigid analytically trivial Anderson $`t`$-motives is equivalent to the category of uniformizable abelian $`t`$-modules defined over $`\overline{k}`$, as in . For a given Anderson $`t`$-motive $`𝖬`$ and associated $`t`$-module $`E`$, there is an explicit connection
$$\text{periods of }E\overline{k}\text{-linear combinations of entries of }\mathrm{\Psi }(\theta )^1.$$
The details of this relationship will be the subject of a future paper with Anderson, but examples are already seen in §3.3 for the Carlitz motive (see also S. K. Sinha \[28, §5.2\] for examples involving special values of the function field $`\mathrm{\Gamma }`$-function).
#### 1.1.5. Remarks on $`t`$-motive terminology
G. Anderson introduced $`t`$-motives in . Later in dual $`t`$-motives, which had several technical advantages, were introduced. The algebraic properties of these two types of $`t`$-motives are essentially the same, and the two categories are anti-equivalent to each other. In this paper we will follow the dual $`t`$-motive point of view only, and throughout we refer to them as Anderson $`t`$-motives. In the following paragraph we discuss a third type of $`t`$-motive, defined properly in §3.4, which are our primary objects of study.
#### 1.1.6. Tannakian category of $`t`$-motives
In §3.4 we show that the category of rigid analytically trivial Anderson $`t`$-motives up to isogeny embeds as a full subcategory of a neutral Tannakian category $`𝒯`$ over $`𝔽_q(t)`$. Objects in $`𝒯`$ are called simply $`t`$-motives, and throughout the paper the term “$`t`$-motive” will refer exclusively to an object in $`𝒯`$. In particular, from this standpoint all $`t`$-motives are rigid analytically trivial. Also objects in $`𝒯`$ do not necessarily come from pure Anderson $`t`$-motives in the sense of , and so $`𝒯`$ is a mixed category.
By Tannakian duality, for each object $`M`$ in $`𝒯`$, the Tannakian subcategory $`𝒯_M`$ generated by $`M`$ satisfies an equivalence of categories
$$𝒯_M\mathrm{𝐑𝐞𝐩}(\mathrm{\Gamma }_M,𝔽_q(t)),$$
where $`\mathrm{𝐑𝐞𝐩}(\mathrm{\Gamma }_M,𝔽_q(t))`$ is the category of finite dimensional representations over $`𝔽_q(t)`$ of some algebraic subgroup $`\mathrm{\Gamma }_M\mathrm{GL}_r`$ defined over $`𝔽_q(t)`$ (see §3.5). The group $`\mathrm{\Gamma }_M`$ is called the Galois group of $`M`$.
It should be noted that R. Pink has defined a category $``$ of mixed Hodge structures for function fields that is a neutral Tannakian category over $`𝔽_q(t)`$. He showed that the category of rigid analytically trivial Anderson $`t`$-motives that are also “mixed” embeds as a full subcategory of $``$. It would be interesting to investigate the relationships among Pink’s Hodge structures, the $`t`$-motives defined in this paper, and their associated Galois groups. In the end our category of $`t`$-motives is best suited for our transcendence applications, so we do not pursue further here the connections with Pink’s work. See also D. Goss for additional comparisons between $`t`$-motives and motives over $``$.
The following is the main theorem of this paper (restated later as Theorem 5.2.2).
###### Theorem 1.1.7.
Let $`M`$ be a $`t`$-motive, and let $`\mathrm{\Gamma }_M`$ be its Galois group. Suppose that $`\mathrm{\Phi }\mathrm{GL}_r(\overline{k}(t))\mathrm{Mat}_r(\overline{k}[t])`$ represents multiplication by $`𝛔`$ on $`M`$ and that $`det\mathrm{\Phi }=c(t\theta )^s`$, $`c\overline{k}^\times `$. Let $`\mathrm{\Psi }`$ be a rigid analytic trivialization of $`\mathrm{\Phi }`$ in $`\mathrm{GL}_r(𝕋)\mathrm{Mat}_r(𝔼)`$. Finally, let $`L`$ be the subfield of $`\overline{k_{\mathrm{}}}`$ generated over $`\overline{k}`$ by the entries of $`\mathrm{\Psi }(\theta )`$. Then
$$\text{tr. deg}_{\overline{k}}L=dim\mathrm{\Gamma }_M.$$
#### 1.1.8. Grothendieck’s conjecture
In light of §1.1.4, the statement of Theorem 1.1.7 can be thought of as a function field version of Grothendieck’s conjecture on periods of algebraic varieties. For an abelian variety $`A`$ over $`\overline{}`$ of dimension $`d`$, let $`P`$ be the period matrix of $`A`$ that represents an isomorphism between $`H^1(A(),)_{}`$ and $`H_{\mathrm{DR}}^1(A/)`$, with basis defined over $`\overline{}`$. Grothendieck’s conjecture is that
$$\text{tr. deg}_\overline{}\overline{}(P)=dim\mathrm{MT}(A),$$
where $`\mathrm{MT}(A)`$ is the Mumford-Tate group of $`A`$ and is an algebraic subgroup of $`\mathrm{GL}_{2d}\times 𝔾_\mathrm{m}`$ over $``$. P. Deligne \[11, Cor. I.6.4\] has proved that the dimension of $`\mathrm{MT}(A)`$ is an upper bound for the transcendence degree. Conjecturally the Mumford-Tate group is isomorphic to the motivic Galois group of the motive $`h_1(A)(1)`$ over $``$. More generally Grothendieck’s period conjecture states that if $`X`$ is a smooth variety over $`\overline{}`$, then
$$\text{tr. deg}_\overline{}\overline{}(P(X))=dim\mathrm{\Gamma }_X^{\text{mot}},$$
where $`P(X)`$ is the period matrix of $`X`$ and $`\mathrm{\Gamma }_X^{\text{mot}}`$ is the motivic Galois group of $`X`$ over $``$. It should be pointed out that by work of C. Bertolin many standard transcendence conjectures over $`\overline{}`$, such as Schanuel’s conjecture, follow from expanded versions of Grothendieck’s period conjecture.
### 1.2. Algebraic independence of Carlitz logarithms
One application of Theorem 1.1.7 is a characterization of algebraic relations over $`\overline{k}`$ of Carlitz logarithms of algebraic numbers.
#### 1.2.1. Carlitz exponential
The Carlitz exponential is the power series
$$\mathrm{exp}_C(z):=z+\underset{i=1}{\overset{\mathrm{}}{}}\frac{z^{q^i}}{(\theta ^{q^i}\theta )(\theta ^{q^i}\theta ^q)\mathrm{}(\theta ^{q^i}\theta ^{q^{i1}})}.$$
As is well known (see \[15, Ch. 3\], \[31, §2.5\]), the function defined by $`\mathrm{exp}_C`$ converges everywhere on $`𝕂`$, is $`𝔽_q`$-linear, and has kernel $`𝔽_q[\theta ]\stackrel{~}{\pi }`$, where
$$\stackrel{~}{\pi }:=\theta \sqrt[q1]{\theta }\underset{i=1}{\overset{\mathrm{}}{}}\left(1\theta ^{1q^i}\right)^1k_{\mathrm{}}(\sqrt[q1]{\theta })^\times .$$
The Carlitz exponential also satisfies the functional equation
$$\mathrm{exp}_C(\theta z)=\theta \mathrm{exp}_C(z)+\mathrm{exp}_C(z)^q,z𝕂.$$
Moreover, this functional equation induces an exact sequence of $`𝔽_q[t]`$-modules,
$$0𝔽_q[\theta ]\stackrel{~}{\pi }𝕂(𝕂)0,$$
where $`(𝕂)`$ is the $`𝔽_q[t]`$-module of $`𝕂`$-valued points on the Carlitz module $``$ (see §3.4.4) and where $`t`$ acts by multiplication by $`\theta `$ on the first two terms. The number $`\stackrel{~}{\pi }`$ is called the Carlitz period.
#### 1.2.2. Carlitz logarithm
The Carlitz logarithm is the inverse of $`\mathrm{exp}_C(z)`$,
$$\mathrm{log}_C(z):=z+\underset{i=1}{\overset{\mathrm{}}{}}\frac{z^{q^i}}{(\theta \theta ^q)(\theta \theta ^{q^2})\mathrm{}(\theta \theta ^{q^i})},$$
which as a function on $`𝕂`$ converges for all $`z𝕂`$ with $`\left|z\right|_{\mathrm{}}<\left|\theta \right|_{\mathrm{}}^{q/(q1)}`$. The Carlitz logarithm is $`𝔽_q`$-linear and satisfies the functional equation
$$\theta \mathrm{log}_C(z)=\mathrm{log}_C(\theta z)+\mathrm{log}_C(z^q),$$
for all $`z𝕂`$ where all three terms converge.
#### 1.2.3. Linear forms in Carlitz logarithms
We recall a theorem of J. Yu. Suppose $`\lambda _1,\mathrm{},\lambda _r𝕂`$ satisfy $`\mathrm{exp}_C(\lambda _i)\overline{k}`$ for each $`i=1,\mathrm{},r`$. As in the previous section there are many potential $`k`$-linear relations among $`\lambda _1,\mathrm{},\lambda _r`$. However, Yu proved that these are the only possible linear relations over $`\overline{k}`$ in the following function field analogue of Baker’s theorem on linear forms in logarithms.
###### Theorem 1.2.4 (Yu \[33, Thm. 4.3\]).
Suppose $`\lambda _1,\mathrm{},\lambda _r𝕂`$ satisfy $`\mathrm{exp}_C(\lambda _i)\overline{k}`$ for $`i=1,\mathrm{},r`$. If $`\lambda _1,\mathrm{},\lambda _r`$ are linearly independent over $`k`$, then the numbers $`1,\lambda _1,\mathrm{},\lambda _r`$ are linearly independent over $`\overline{k}`$.
Yu’s result is an application of his far reaching Theorem of the Sub-$`t`$-module \[33, Thm. 0.1\], which characterizes all $`\overline{k}`$-linear relations among logarithms of points in $`\overline{k}`$ on general $`t`$-modules. Transcendence results about the Carlitz periods and Carlitz logarithms go back to Carlitz and Wade in the 1940’s. For detailed accounts of the history of transcendence results for Drinfeld modules, including Yu’s theorem, see W. D. Brownawell and D. S. Thakur \[31, Ch. 10\].
#### 1.2.5. Algebraic independence of Carlitz logarithms
In characteristic $`0`$, Baker’s theorem on linear forms in natural logarithms of algebraic numbers is best known. In the situation of Carlitz logarithms we use Theorem 1.1.7 to prove the following theorem (restated later as Theorem 6.4.2).
###### Theorem 1.2.6.
Let $`\lambda _1,\mathrm{},\lambda _r𝕂`$ satisfy $`\mathrm{exp}_C(\lambda _i)\overline{k}`$ for each $`i=1,\mathrm{},r`$. If $`\lambda _1,\mathrm{},\lambda _r`$ are linearly independent over $`k`$, then they are algebraically independent over $`\overline{k}`$.
It should be noted that, using Mahler’s method, L. Denis has proved the special case of this theorem where $`\lambda _1,\mathrm{},\lambda _r`$ are restricted to values of $`\mathrm{log}_C`$ on elements of $`𝔽_q(\theta ^{1/e})`$, $`e1`$, of degree in $`\theta `$ less than $`q/(q1)`$.
### 1.3. Methods of proof
#### 1.3.1. $`\sigma `$-semilinear difference equations
The category of $`t`$-motives is a certain full subcategory in the category of left $`\overline{k}(t)[𝝈,𝝈^1]`$-modules which are finite dimensional as $`\overline{k}(t)`$-vector spaces. To every $`t`$-motive $`M`$ one can associate a matrix $`\mathrm{\Phi }\mathrm{GL}_r(\overline{k}(t))`$ representing multiplication by $`\sigma `$ and a rigid analytic trivialization $`\mathrm{\Psi }\mathrm{GL}_r(𝕃)`$ so that $`\mathrm{\Psi }^{(1)}=\mathrm{\Phi }\mathrm{\Psi }`$. Here recall that $`𝕃`$ is the fraction field of the Tate algebra $`𝕋`$. Thus the columns of $`\mathrm{\Psi }`$ satisfy a system of $`\sigma `$-semilinear difference equations in the sense of , where $`\sigma =(ff^{(1)}):𝕃\stackrel{}{}𝕃`$, and we develop the theory of such equations in this context in §4. In spirit this theory is close to the Galois theory of differential equations and difference equations in characteristic $`0`$ , , , , , , .
In §4 we develop the Picard-Vessiot theory for certain kinds of difference equations for $`\sigma `$ and construct their difference Galois groups (see Theorem 4.2.11). However, careful attention must be paid to the fact that the fixed field of $`\sigma `$ in $`\overline{k}(t)`$ is $`𝔽_q(t)`$. The Galois theory of difference equations developed by M. van der Put and M. F. Singer is quite useful here, but it does not completely apply because they fundamentally use that the field of fixed elements under the difference automorphism is algebraically closed. On the one hand, because the fixed field of $`\sigma `$ in $`𝕃`$ is also $`𝔽_q(t)`$, the Galois groups we construct are themselves defined over $`𝔽_q(t)`$. However, that $`𝔽_q(t)`$ is not algebraically closed nor even perfect presents several difficulties because in general the $`𝔽_q(t)`$-valued points of the Galois group need not be dense and the group itself need not be a priori smooth.
#### 1.3.2. $`t`$-motives and difference Galois groups
Given a $`t`$-motive $`M`$ of dimension $`r`$ over $`\overline{k}(t)`$, the difference Galois group $`\mathrm{\Gamma }`$ is a subgroup of $`\mathrm{GL}_r`$ over $`𝔽_q(t)`$. Let $`\mathrm{\Sigma }`$ be the $`\overline{k}(t)`$-subalgebra of $`𝕃`$ generated by the entries of $`\mathrm{\Psi }`$ and $`det(\mathrm{\Psi })^1`$, and let $`\mathrm{\Lambda }`$ be its fraction field. The field $`𝕃`$ is naturally a left $`\overline{k}(t)[𝝈,𝝈^1]`$-module via the automorphism $`\sigma `$, and $`\mathrm{\Sigma }`$ and $`\mathrm{\Lambda }`$ are both $`\sigma `$-invariant. Then
$$\mathrm{\Gamma }(𝔽_q(t))\mathrm{Aut}_\sigma (\mathrm{\Sigma }/\overline{k}(t)),$$
where the right-hand side is the group of automorphisms of $`\mathrm{\Sigma }`$ over $`\overline{k}(t)`$ that commute with $`\sigma `$. Moreover, this identification is compatible with base extensions of $`𝔽_q(t)`$ (see §4.4.14.4.3).
We work out an explicit description of $`\mathrm{\Gamma }(\overline{𝔽_q(t)})`$ in §4.4, and, using crucially that $`𝕃`$ is a separable extension of $`\overline{k}(t)`$ and that $`\overline{k}(t)`$ is algebraically closed in $`\mathrm{\Lambda }`$, we show that $`\mathrm{\Gamma }`$ has the following properties:
* $`\mathrm{\Gamma }`$ is smooth over $`𝔽_q(t)`$ (Theorem 4.3.1(b));
* $`dim\mathrm{\Gamma }=\text{tr. deg}_{\overline{k}(t)}\mathrm{\Lambda }`$, (Theorem 4.3.1(c));
* The elements of $`\mathrm{\Lambda }`$ fixed by $`\mathrm{\Gamma }(\overline{𝔽_q(t)})`$ are precisely $`\overline{k}(t)`$ (Theorem 4.4.6).
These properties are essential for proving in Theorem 4.5.10 that
$$\mathrm{\Gamma }\mathrm{\Gamma }_M,$$
where $`\mathrm{\Gamma }_M`$ is the Galois group associated to $`M`$ by Tannakian duality.
#### 1.3.3. The proof of Theorem 1.1.7
The primary vehicle for proving this theorem is a $`\overline{k}`$-linear independence criterion from \[2, Thm. 3.1.1\]. It is stated here in Theorem 5.1.1. We apply this criterion to the rigid analytic trivializations of tensor powers of $`M`$ so as to compare the dimensions of the $`\overline{k}`$-span of monomials of the entries of $`\mathrm{\Psi }(\theta )`$ of a given degree and the $`\overline{k}(t)`$-span of monomials in the entries of $`\mathrm{\Psi }`$. Ultimately we show that
$$\text{tr. deg}_{\overline{k}}L=\text{tr. deg}_{\overline{k}(t)}\mathrm{\Lambda },$$
the latter of which is the same as the dimension of $`\mathrm{\Gamma }_M`$.
#### 1.3.4. Carlitz logarithms
For $`\alpha _1,\mathrm{},\alpha _r\overline{k}^\times `$ with $`\left|\alpha _i\right|_{\mathrm{}}<\left|\theta \right|_{\mathrm{}}^{q/(q1)}`$ for $`i=1,\mathrm{},r`$, we define a $`t`$-motive $`X`$ so that the field generated over $`\overline{k}`$ by the entries of its rigid analytic trivialization $`\mathrm{\Psi }`$ evaluated at $`t=\theta `$ is precisely
$$L=\overline{k}(\mathrm{\Psi }(\theta ))=\overline{k}(\stackrel{~}{\pi },\mathrm{log}_C(\alpha _1),\mathrm{},\mathrm{log}_C(\alpha _r)).$$
Moreover, we show that arbitrary logarithms are $`k`$-linear combinations of logarithms of this form in a precise way. We determine a set of defining equations of the Galois group $`\mathrm{\Gamma }_X`$ of $`X`$ in Theorem 6.3.2 each of which is a linear polynomial over $`\overline{k}(t)`$. These linear relations each produce a $`k`$-linear relation on the logarithms and $`\stackrel{~}{\pi }`$. We then use Theorem 1.1.7 to show that all algebraic relations must arise from these relations.
### 1.4. Acknowledgements
The author thanks D. Brownawell, L. Denis, C.-Y. Chang, D. Goss, L.-C. Hsia, M. van der Put, D. Thakur, and J. Yu for many helpful discussions on the contents of this paper. He further thanks the National Center for Theoretical Sciences in Hsinchu, Taiwan, where many of the results in this paper were proved. The author especially thanks G. Anderson and N. Ramachandran for their indispensable advice throughout this project. Finally the author thanks the referee for several useful suggestions.
## 2. Notation and preliminaries
### 2.1. Table of symbols
### 2.2. Preliminaries
#### 2.2.1. Norms
We let $`||_{\mathrm{}}`$ denote a fixed $`\mathrm{}`$-adic norm on $`𝕂`$. For a matrix $`E\mathrm{Mat}_{r\times s}(𝕂)`$, we set $`\left|E\right|_{\mathrm{}}=sup\left|E_{ij}\right|_{\mathrm{}}`$. For matrices $`E`$ and $`F`$, we observe that $`\left|E+F\right|_{\mathrm{}}\mathrm{max}(\left|E\right|_{\mathrm{}},\left|F\right|_{\mathrm{}})`$ and $`\left|EF\right|_{\mathrm{}}\left|E\right|_{\mathrm{}}\left|F\right|_{\mathrm{}}`$.
#### 2.2.2. Generalized power series
Let $`F`$ be a field of characteristic $`p`$. For a formal series $`f:=_ia_it^i`$ with $`a_iF`$, we let $`\mathrm{Supp}(f):=\{ia_i0\}`$. We let $`Ft`$ be the set of such series for which $`\mathrm{Supp}(f)`$ is a well-ordered subset of $``$. This condition implies that $`Ft`$ is a field under the natural addition and multiplication of these series so that $`t^it^j=t^{i+j}`$ (see P. Ribenboim \[27, §2\]). If $`F`$ is algebraically closed, then $`Ft`$ is algebraically closed \[27, §5\]. If $`F`$ is a perfect field, then $`Ft`$ is also perfect.
It should be noted that, when $`F`$ is algebraically closed, $`Ft`$ is not the algebraic closure of the Laurent series field $`F((t))`$. For an explicit description of the field $`\overline{F((t))}Ft`$, the reader is directed to K. Kedlaya .
By considering the inclusions
$$𝔽_q(t)\overline{k}(t)𝕂((t))𝕂t,$$
we fix once and for all the inclusions of algebraically closed fields
$$\overline{𝔽_q(t)}\overline{k(t)}\overline{𝕂((t))}𝕂t.$$
#### 2.2.3. Entire functions
A power series $`f=_{i=0}^{\mathrm{}}a_it^i𝕂[[t]]`$ that satisfies
$$\underset{i\mathrm{}}{lim}\sqrt[i]{\left|a_i\right|_{\mathrm{}}}=0$$
and
$$[k_{\mathrm{}}(a_0,a_1,a_2,\mathrm{}):k_{\mathrm{}}]<\mathrm{},$$
is an *entire power series*. As a function of $`t`$, such a power series $`f`$ converges on all of $`𝕂`$, and, when restricted to $`\overline{k_{\mathrm{}}}`$, $`f`$ takes values in $`\overline{k_{\mathrm{}}}`$. The ring of entire power series is denoted $`𝔼`$.
#### 2.2.4. Restricted Laurent series
A power series $`_{i=0}^{\mathrm{}}a_it^i𝕂[[t]]`$ that satisfies
$$\underset{i\mathrm{}}{lim}\left|a_i\right|_{\mathrm{}}=0,$$
is called a *restricted power series*. As functions of $`t`$, these power series converge on the closed unit disk in $`𝕂`$. The restricted power series form a subring $`𝕋=𝕂\{t\}`$ of $`𝕂[[t]]`$, and $`𝔼`$ is a subring of $`𝕋`$. The fraction field of $`𝕋`$, denoted $`𝕃`$, is the field of *restricted Laurent series*.
Now at each point $`a𝕂`$ with $`\left|a\right|_{\mathrm{}}1`$, a function $`f𝕃`$ has a well-defined order of vanishing $`\mathrm{ord}_a(f)`$, and for all but finitely many $`\left|a\right|_{\mathrm{}}1`$, we have $`\mathrm{ord}_a(f)=0`$. Also each $`f𝕃`$ has a unique factorization
(2.2.4.1)
$$f=\lambda \left[\underset{|a|_{\mathrm{}}1}{}(ta)^{\mathrm{ord}_a(f)}\right]\left[1+\underset{i=1}{\overset{\mathrm{}}{}}b_it^i\right],$$
where $`0\lambda 𝕂`$, $`sup\left|b_i\right|_{\mathrm{}}<1`$, and $`\left|b_i\right|_{\mathrm{}}0`$ (see \[13, Cor. 2.2.4\]). The series $`1+b_it^i`$ is a unit in $`𝕋`$, and it follows that $`𝕋`$ is a principal ideal domain with maximal ideals generated by each $`ta`$, $`\left|a\right|_{\mathrm{}}1`$ (see \[13, Thm. 2.2.9\]).
For $`f=_{i=0}^{\mathrm{}}a_it^i𝕋`$, we define its norm $`f`$ to be
$$f:=\underset{i}{sup}\left|a_i\right|_{\mathrm{}}=\underset{i}{\mathrm{max}}\left|a_i\right|_{\mathrm{}}.$$
If $`f𝕋`$ is written as in (2.2.4.1), then $`f=\left|\lambda \right|_{\mathrm{}}`$. The norm $``$ is a complete ultrametric norm on $`𝕋`$ and satisfies
$`cf`$ $`=\left|c\right|_{\mathrm{}}f,`$ $`c𝕂,f𝕋,`$
$`fg`$ $`=fg,`$ $`f,g𝕋.`$
#### 2.2.5. Twisting
We define an automorphism $`\sigma :𝕂t𝕂t`$ by setting
$$\sigma \left(\underset{i}{}a_it^i\right):=\underset{i}{}a_i^{1/q}t^i.$$
If $`f𝕂t`$ and $`n`$, the *$`n`$-fold twist* of $`f`$ is defined to be
$$f^{(n)}:=\sigma ^n(f).$$
The automorphism $`\sigma `$ of $`𝕂t`$ induces automorphisms of several subrings, notably $`\overline{k}[t]`$, $`\overline{k}(t)`$, $`𝔼`$, $`𝕋`$, $`𝕃`$, $`𝕂[[t]]`$, $`𝕂((t))`$. Moreover, $`\sigma `$ also leaves $`\overline{𝔽_q(t)}`$, $`\overline{k(t)}`$, $`\overline{𝕃}`$, and $`\overline{𝕂((t))}`$ invariant.
If $`F`$ is a subring of $`𝕂t`$ that is invariant under $`\sigma `$, we set
$$F^\sigma :=\{fF\sigma (f)=f\}$$
to be the elements of $`F`$ fixed by $`\sigma `$. It is clear that $`F^{\sigma ^n}`$ is a subring of $`F`$ and that $`F^{\sigma ^m}F^{\sigma ^n}`$ if $`mn`$. For example,
$`𝕂t^\sigma `$ $`=𝔽_qt,`$ $`\overline{k(t)}^\sigma `$ $`=\overline{𝔽_q(t)}^\sigma =\overline{𝔽_q(t)}𝔽_qt,`$
$`𝕂((t))^\sigma `$ $`=𝔽_q((t)),`$ $`\overline{k}(t)^\sigma `$ $`=𝔽_q(t).`$
The only item that requires any explanation here is the description of $`\overline{k(t)}^\sigma `$. For $`\alpha \overline{k(t)}^\sigma `$, let $`x^m+b_{m1}x^{m1}+\mathrm{}+b_0\overline{k}(t)[x]`$ be the minimal polynomial of $`\alpha `$ over $`\overline{k}(t)`$. Since $`\sigma (\alpha )=\alpha `$, we have that $`\alpha `$ is also a root of $`x^m+\sigma (b_{m1})x^{m1}+\mathrm{}+\sigma (b_0)`$. Taking the difference of these two relations, we see that $`\sigma (b_i)=b_i`$ for each $`i`$, and so the minimal polynomial of $`\alpha `$ has coefficients in $`\overline{k}(t)^\sigma =𝔽_q(t)`$.
If $`F`$ is a matrix with entries in $`𝕂t`$, then $`\sigma ^n(F):=F^{(n)}`$ is defined by the rule $`\left(F^{(n)}\right)_{ij}:=\left(F_{ij}^{(n)}\right)`$. If $`F\mathrm{Mat}_{r\times s}(𝕃)`$, we set $`F:=\mathrm{max}_{i,j}F_{ij}`$, in which case $`F^{(n)}=F^{q^n}`$.
###### Lemma 2.2.6.
For any $`\alpha 𝕂`$, there is a positive integer $`s`$ so that with respect to $`||_{\mathrm{}}`$ on $`𝕂`$,
$$\underset{n\mathrm{}}{lim}\alpha ^{(ns)}=\{\begin{array}{cc}0\hfill & \text{if }\left|\alpha \right|_{\mathrm{}}<1\text{,}\hfill \\ c\overline{𝔽}_q^\times \hfill & \text{if }\left|\alpha \right|_{\mathrm{}}=1\text{,}\hfill \\ \mathrm{}\hfill & \text{if }\left|\alpha \right|_{\mathrm{}}>1\text{.}\hfill \end{array}$$
###### Proof.
If $`\left|\alpha \right|_{\mathrm{}}1`$, then the result is clear. Otherwise, there is a unique $`c\overline{𝔽}_q^\times `$ so that $`\left|\alpha c\right|_{\mathrm{}}<1`$. (See \[28, Lem. 2.4.4\].) Then $`c𝔽_{q^s}`$ for some $`s1`$, and the result follows. ∎
###### Lemma 2.2.7.
For any $`f𝕋`$ with $`f1`$, there is a positive integer $`s`$ so that with respect to $``$ on $`𝕋`$,
$$\underset{n\mathrm{}}{lim}f^{(ns)}\overline{𝔽}_q[t].$$
Also $`f=1`$ if and only if $`lim_n\mathrm{}f^{(ns)}0`$.
###### Proof.
We use the factorization of $`f`$ in (2.2.4.1). For each $`a`$, with $`\left|a\right|_{\mathrm{}}1`$, if $`\mathrm{ord}_a(f)0`$, then as in Lemma 2.2.6 choose $`s_a1`$ and $`c_a\overline{𝔽}_q`$ so that $`lim_n\mathrm{}a^{(ns_a)}=c_a`$. Likewise, since $`f1`$, we have $`\left|\lambda \right|_{\mathrm{}}1`$, and so we can choose $`s_\lambda 1`$ and $`c_\lambda \overline{𝔽}_q`$ with $`\lambda ^{(ns_\lambda )}c_\lambda `$. Then we let $`s`$ be the least common multiple of all the $`s_a`$’s and $`s_\lambda `$. From (2.2.4.1), with respect to $``$,
$$\underset{n\mathrm{}}{lim}\left[1+\underset{i=1}{\overset{\mathrm{}}{}}b_it^i\right]^{(ns)}=1,$$
since $`sup\left|b_i\right|_{\mathrm{}}<1`$. Therefore,
$$\underset{n\mathrm{}}{lim}f^{(ns)}=\underset{n\mathrm{}}{lim}\lambda ^{(ns)}\underset{\left|a\right|_{\mathrm{}}1}{}\left(ta^{(ns)}\right)^{\mathrm{ord}_a(f)}=c_\lambda \underset{\left|a\right|_{\mathrm{}}1}{}(tc_a)^{\mathrm{ord}_a(f)},$$
which is in $`\overline{𝔽}_q[t]`$. Furthermore, $`f=1`$ if and only if $`\left|\lambda \right|_{\mathrm{}}=1`$, which holds if and only if $`c_\lambda 0`$. Thus $`f=1`$ if and only if $`lim_n\mathrm{}f^{(ns)}0`$. ∎
## 3. $`t`$-motives and Tannakian categories
Here we will define a category $`𝒯`$ of $`t`$-motives that is a neutral Tannakian category over $`𝔽_q(t)`$. For all definitions of tensor categories and Tannakian categories, we follow Deligne and J. S. Milne \[11, §II\]. Other useful references include , , \[26, App. B\].
As mentioned in §1.1.6, Tannakian categories for $`t`$-motives have been considered previously by Pink , though through a different construction. Parts of the theory of $`t`$-motives defined below have been considered by Y. Taguchi and A. Tamagawa in their study of the Tate conjecture for $`t`$-modules. Also our theory has similarities with the theory of $`\sigma `$-bundles defined by U. Hartl and Pink .
### 3.1. The rings $`\overline{k}[t;𝝈]`$ and $`\overline{k}(t)[𝝈,𝝈^1]`$
#### 3.1.1. Definition
The ring $`\overline{k}(t)[𝝈,𝝈^1]`$ is the noncommutative ring of Laurent polynomials in the variable $`𝝈`$ with coefficients in $`\overline{k}(t)`$, subject to the relation
$$𝝈f=\sigma (f)𝝈=f^{(1)}𝝈$$
for all $`f\overline{k}(t)`$. Thus every element of $`\overline{k}(t)[𝝈,𝝈^1]`$ has the form $`_{i=m}^mf_i𝝈^i`$, where $`f_i\overline{k}(t)`$.
#### 3.1.2. Ring-theoretic properties
The polynomials in $`𝝈`$ with coefficients in $`\overline{k}[t]`$ comprise the subring $`\overline{k}[t;𝝈]`$ of $`\overline{k}(t)[𝝈,𝝈^1]`$. The ring $`\overline{k}[𝝈]`$ is the subring of polynomials with coefficients in $`\overline{k}`$. Both $`\overline{k}[t;𝝈]`$ and $`\overline{k}(t)[𝝈,𝝈^1]`$ are domains. The center of $`\overline{k}[t;𝝈]`$ is $`𝔽_q[t]`$, and the center of $`\overline{k}(t)[𝝈,𝝈^1]`$ is $`𝔽_q(t)`$. The fundamental properties of the ring $`\overline{k}[t;𝝈]`$ are covered in \[2, §4\].
### 3.2. Pre-$`t`$-motives
Here we define the category $`𝒫`$ of pre-$`t`$-motives and explore its basic properties. In particular we show in Theorem 3.2.13 that $`𝒫`$ is a rigid abelian $`𝔽_q(t)`$-linear tensor category.
#### 3.2.1. The category $`𝒫`$
We let $`𝒫`$ be the category of left $`\overline{k}(t)[𝝈,𝝈^1]`$-modules that are finite dimensional over $`\overline{k}(t)`$. Morphisms in $`𝒫`$ are left $`\overline{k}(t)[𝝈,𝝈^1]`$-module homomorphisms. We call $`𝒫`$ the category of *pre-$`t`$-motives*, though it is worth noting that $`𝒫`$ is a category of difference modules with respect to the automorphism $`\sigma :\overline{k}(t)\overline{k}(t)`$ in the sense of .
#### 3.2.2. Preliminary properties of $`𝒫`$
The category of pre-$`t`$-motives is an abelian category. For two objects $`P`$ and $`Q`$ in $`𝒫`$, it follows that $`\mathrm{Hom}_𝒫(P,Q)`$ is an $`𝔽_q(t)`$-vector space. A straightforward adaptation of the proof of \[1, Thm. 2\] shows that the map
$$\mathrm{Hom}_𝒫(P,Q)_{𝔽_q(t)}\overline{k}(t)\mathrm{Hom}_{\overline{k}(t)}(P,Q)$$
is injective. Thus $`\mathrm{Hom}_𝒫(P,Q)`$ is a finite dimensional $`𝔽_q(t)`$-vector space.
#### 3.2.3. Representations of pre-$`t`$-motives
Given a $`\overline{k}(t)`$-vector space $`P`$ and $`p_1,\mathrm{},p_rP`$, we call the vector
$$𝐩=\left[\begin{array}{c}p_1\\ \mathrm{}\\ p_r\end{array}\right]\mathrm{Mat}_{r\times 1}(P)$$
a *basis for $`P`$* if $`p_1,\mathrm{},p_r`$ form a $`\overline{k}(t)`$-basis for $`P`$. If $`P`$ is a pre-$`t`$-motive, then there is a unique matrix $`\mathrm{\Phi }=\mathrm{\Phi }_𝐩\mathrm{GL}_r(\overline{k}(t))`$ such that
$$𝝈𝐩=\mathrm{\Phi }𝐩.$$
We say that *$`\mathrm{\Phi }`$ represents multiplication by $`𝛔`$ on $`P`$*. Moreover, the matrix $`\mathrm{\Phi }\mathrm{GL}_r(\overline{k}(t))`$ uniquely determines the left $`\overline{k}(t)[𝝈,𝝈^1]`$-module structure on $`P`$ with respect to $`𝐩`$.
Now suppose that $`\varphi :PQ`$ is a morphism of pre-$`t`$-motives and that $`𝐩\mathrm{Mat}_{r\times 1}(P)`$ and $`𝐪\mathrm{Mat}_{s\times 1}(Q)`$ are bases for $`P`$ and $`Q`$ respectively. If $`B\mathrm{Mat}_{r\times s}(\overline{k}(t))`$ represents $`\varphi `$ as a map of $`\overline{k}(t)`$-vector spaces such that
$$\varphi (𝐟𝐩)=𝐟B𝐪,𝐟\mathrm{Mat}_{1\times r}(\overline{k}(t)),$$
then
$$B^{(1)}\mathrm{\Phi }_𝐪=\mathrm{\Phi }_𝐩B.$$
In particular, if $`𝐪`$ is simply another basis of $`P`$, and $`B\mathrm{GL}_r(\overline{k}(t))`$ is the change of basis matrix, then $`\mathrm{\Phi }_𝐩=B^{(1)}\mathrm{\Phi }_𝐪B^1`$.
#### 3.2.4. Tensor products of pre-$`t`$-motives
Let $`P`$ and $`Q`$ be pre-$`t`$-motives. Then the $`\overline{k}(t)`$-vector space $`P_{\overline{k}(t)}Q`$ is made into a $`\overline{k}(t)[𝝈,𝝈^1]`$-module by defining
$$𝝈(mn):=(𝝈m)(𝝈n).$$
It is clear that then multiplication by $`𝝈`$ is bijective on $`P_{\overline{k}(t)}Q`$ and that $`P_{\overline{k}(t)}Q`$ is a pre-$`t`$-motive. Likewise we define arbitrary finite tensor products of pre-$`t`$-motives with diagonal $`𝝈`$-action. For a fixed pre-$`t`$-motive $`P`$ and $`n1`$, we set $`P^n:=_{i=1}^nP`$ to be the $`n`$-th tensor power of $`P`$.
#### 3.2.5. Representations of tensor products
Let $`𝐩=[p_1,\mathrm{},p_r]^{\text{tr}}`$ and $`𝐪=[q_1,\mathrm{},q_s]^{\text{tr}}`$ be $`\overline{k}(t)`$-bases for pre-$`t`$-motives $`P`$ and $`Q`$ respectively. Then, with respect to the basis
$$𝐩𝐪:=[p_1q_1,p_1q_2,\mathrm{},p_rq_s]^{\text{tr}},$$
on $`PQ`$, the Kronecker product, $`\mathrm{\Phi }_{𝐩𝐪}=\mathrm{\Phi }_𝐩\mathrm{\Phi }_𝐪`$, represents multiplication by $`𝝈`$ on $`PQ`$. Similarly these conventions extend to arbitrary finite tensor products of pre-$`t`$-motives.
#### 3.2.6. The Carlitz motive
We define the *Carlitz motive* to be the pre-$`t`$-motive $`C`$ whose underlying $`\overline{k}(t)`$-vector space is $`\overline{k}(t)`$ itself and on which $`𝝈`$ acts by
$$𝝈f:=(t\theta )f^{(1)},fC.$$
For $`n1`$, the underlying $`\overline{k}(t)`$-vector space of $`C^n`$ is also $`\overline{k}(t)`$, and multiplication by $`𝝈`$ on $`C^n`$ is given by
$$𝝈f=(t\theta )^nf^{(1)},fC^n.$$
See also .
#### 3.2.7. Internal Hom
Let $`P`$ and $`Q`$ be pre-$`t`$-motives, and set
$$R:=\mathrm{Hom}_{\overline{k}(t)}(P,Q).$$
Then $`R`$ is a $`\overline{k}(t)`$-vector space. We define a $`\overline{k}(t)[𝝈,𝝈^1]`$-module structure on $`R`$ by setting
$$𝝈\rho :=𝝈\rho 𝝈^1,\rho R.$$
It is straightforward to check that $`𝝈\rho :PQ`$ is $`\overline{k}(t)`$-linear, and so $`𝝈:RR`$, and that this action of $`𝝈`$ extends naturally to a left $`\overline{k}(t)[𝝈,𝝈^1]`$-module structure on $`R`$. We write $`\mathrm{Hom}(P,Q)`$ for the $`\overline{k}(t)[𝝈,𝝈^1]`$-module $`R`$ just defined. It is also a pre-$`t`$-motive.
#### 3.2.8. Identity object
Let $`\mathrm{𝟏}:=\overline{k}(t)`$ together with a $`𝝈`$-action defined by
$$𝝈f=\sigma (f)=f^{(1)},f\mathrm{𝟏}.$$
Then $`\mathrm{𝟏}`$ is a pre-$`t`$-motive. Moreover, for any pre-$`t`$-motive $`P`$, the natural isomorphisms, $`P\mathrm{𝟏}\mathrm{𝟏}PP`$, are isomorphisms of pre-$`t`$-motives. Thus $`\mathrm{𝟏}`$ is an identity object with respect to tensor products in $`𝒫`$.
###### Lemma 3.2.9.
$`\mathrm{End}_𝒫(\mathrm{𝟏})𝔽_q(t)`$.
###### Proof.
Suppose $`\varphi :\mathrm{𝟏}\mathrm{𝟏}`$ is a morphism in $`𝒫`$. As a map of $`\overline{k}(t)`$-vector spaces, there is some $`a\overline{k}(t)`$ so that $`\varphi (f)=af`$ for all $`f\overline{k}(t)`$. Since $`\varphi `$ is also $`\overline{k}(t)[𝝈,𝝈^1]`$-linear, we must have $`𝝈a=a𝝈`$, and so $`a`$ is in the center of $`\overline{k}(t)[𝝈,𝝈^1]`$. Thus $`a𝔽_q(t)`$. ∎
#### 3.2.10. Duals
Let $`P`$ be a pre-$`t`$-motive. Then set
$$P^{}:=\mathrm{Hom}(P,\mathrm{𝟏}).$$
The pre-$`t`$-motive $`P^{}`$ is called the *dual of $`P`$*. As a $`\overline{k}(t)`$-vector space, $`P^{}`$ is the dual vector space of $`P`$. If $`𝐩`$ forms a basis for $`P`$, let $`𝐩^{}`$ be the dual basis. We find easily that
$$\mathrm{\Phi }_𝐩^{}=\left(\mathrm{\Phi }_𝐩^1\right)^{\text{tr}}.$$
If $`\varphi :PQ`$ is a morphism of pre-$`t`$-motives, then the dual morphism of $`\overline{k}(t)`$-vector spaces, $`\varphi ^{}:Q^{}P^{}`$, is also $`\overline{k}(t)[𝝈,𝝈^1]`$-linear. These constructions are functorial in $`P`$ and $`Q`$, and thus $`PP^{}:𝒫𝒫`$ defines a contravariant $`𝔽_q(t)`$-linear functor.
#### 3.2.11. Dual of the Carlitz motive
Using the definition of the Carlitz motive in §3.2.6, we see that $`C^{}`$ is isomorphic to $`\overline{k}(t)`$ as a $`\overline{k}(t)`$-vector space and that
$$𝝈f=\frac{1}{t\theta }f^{(1)},fC^{}\text{ (}=\overline{k}(t)\text{).}$$
Furthermore, we see that $`C^{}C\mathrm{𝟏}`$ and that $`C`$ is an invertible object in $`𝒫`$. Thus the functor
$$PPC:𝒫𝒫$$
is an equivalence of categories. We define for $`n`$,
$$C(n):=\{\begin{array}{cc}C^n\hfill & \text{if }n>0\text{,}\hfill \\ \mathrm{𝟏}\hfill & \text{if }n=0\text{,}\hfill \\ (C^{})^n\hfill & \text{if }n<0\text{.}\hfill \end{array}$$
#### 3.2.12. Rigid abelian tensor category
In the language of \[11, §II.1\], it is easily shown that the category of pre-$`t`$-motives is an abelian $`𝔽_q(t)`$-linear tensor category. We omit the details, but we observe that
* each $`\mathrm{Hom}_𝒫(P,Q)`$ is a finite dimensional vector space over $`𝔽_q(t)`$;
* $``$ is compatibly associative and commutative;
* $``$ is $`𝔽_q(t)`$-bilinear;
* $`\mathrm{𝟏}`$ is an identity object with respect to tensor products.
Furthermore, it is straightforward to check that
* the pre-$`t`$-motive $`\mathrm{Hom}(P,Q)`$ defines an internal Hom in $`𝒫`$ that is compatible with tensor products;
* for each pre-$`t`$-motive $`P`$, there is a natural isomorphism $`PP^{}`$.
Therefore, $`𝒫`$ is also rigid. We record this information in the following theorem.
###### Theorem 3.2.13.
The category $`𝒫`$ of pre-$`t`$-motives is a rigid abelian $`𝔽_q(t)`$-linear tensor category.
### 3.3. Rigid analytic triviality
#### 3.3.1. The category $``$
Let $`P`$ be a pre-$`t`$-motive. We set
$$P^{}:=𝕃_{\overline{k}(t)}P,$$
and give $`P^{}`$ a left $`\overline{k}(t)[𝝈,𝝈^1]`$-module structure by setting
$$𝝈(fm):=f^{(1)}𝝈m.$$
Let
$$P^\mathrm{B}:=(P^{})^𝝈=\{\mu P^{}𝝈\mu =\mu \}.$$
Then $`P^\mathrm{B}`$ is an $`𝔽_q(t)`$-vector space, and $`PP^\mathrm{B}`$ is a covariant functor from $`𝒫`$ to the category of $`𝔽_q(t)`$-vector spaces. (The “$`\mathrm{B}`$” in $`P^\mathrm{B}`$ stands for “Betti.”) It is straightforward to check that $`PP^\mathrm{B}`$ is left exact.
We say that $`P`$ is *rigid analytically trivial* if the natural map
$$𝕃_{𝔽_q(t)}P^\mathrm{B}P^{}$$
is an isomorphism. If $`PQ`$ as pre-$`t`$-motives and $`P`$ is rigid analytically trivial, then so is $`Q`$. We let $``$ denote the strictly full subcategory of $`𝒫`$ whose objects are the rigid analytically trivial pre-$`t`$-motives. Clearly the zero object is rigid analytically trivial, and so $``$ is non-empty. We shall see momentarily that $`\mathrm{𝟏}`$ and $`C`$ are also rigid analytically trivial.
###### Lemma 3.3.2.
We have $`𝕃^\sigma =𝔽_q(t)`$.
###### Proof.
By definition, for $`f𝕃^\sigma `$ we have $`f^{(1)}=f`$, and so by (2.2.4.1) the polar divisor $`D`$ of $`f`$ on the closed unit disk in $`𝕂`$ must also satisfy $`D^{(1)}=D`$. Therefore $`D`$ is the divisor of zeros of a polynomial $`c`$ in $`𝔽_q[t]`$. Then $`cf𝕋`$, and $`(cf)^{(1)}=cf`$, from which we have $`cf𝕋𝔽_q[[t]]=𝔽_q[t]`$. ∎
###### Proposition 3.3.3.
The pre-$`t`$-motive $`\mathrm{𝟏}`$ is rigid analytically trivial.
###### Proof.
It is clear that $`\mathrm{𝟏}^{}=𝕃`$ with $`𝝈f:=f^{(1)}`$ for $`f𝕃`$. Therefore, by Lemma 3.3.2, $`\mathrm{𝟏}^\mathrm{B}=𝕃^\sigma =𝔽_q(t)`$. Thus $`𝕃_{𝔽_q(t)}\mathrm{𝟏}^\mathrm{B}\mathrm{𝟏}^{}`$. ∎
#### 3.3.4. The power series $`\mathrm{\Omega }`$
Consider the power series
$$\mathrm{\Omega }=\mathrm{\Omega }(t):=\zeta _\theta ^q\underset{i=1}{\overset{\mathrm{}}{}}\left(1t/\theta ^{(i)}\right)k_{\mathrm{}}(\zeta _\theta )[[t]]𝕂[[t]].$$
It is not difficult to show that $`\mathrm{\Omega }(t)`$ has an infinite radius of convergence, and so $`\mathrm{\Omega }𝔼𝕋`$. Since $`\mathrm{\Omega }`$ has infinitely many zeros in $`𝕂`$, it follows that $`\mathrm{\Omega }\overline{𝕂(t)}`$. Since $`\mathrm{\Omega }`$ has no zeros inside the unit disk, it follows that $`\mathrm{\Omega }𝕋^\times `$. It also satisfies the functional equation
$$\mathrm{\Omega }^{(1)}=(t\theta )\mathrm{\Omega }.$$
The number
$$\stackrel{~}{\pi }=\frac{1}{\mathrm{\Omega }(\theta )}=\theta \zeta _\theta \underset{i=1}{\overset{\mathrm{}}{}}\left(1\theta ^{1q^i}\right)^1k_{\mathrm{}}(\zeta _\theta )$$
is the *Carlitz period*, which figures prominently in our transcendence considerations later on (see also \[3, Cor. 5.2.8\], \[15, §3.2\], \[31, §2.5\]).
###### Lemma 3.3.5.
Suppose $`f𝕃`$ satisfies $`(t\theta )^nf^{(1)}=f`$ for some $`n`$. Then $`f=c/\mathrm{\Omega }^n`$ for some $`c𝔽_q(t)`$.
###### Proof.
Let $`c=f\mathrm{\Omega }^n`$. Then $`c`$ satisfies $`c^{(1)}=c`$, and so by Lemma 3.3.2, $`c𝔽_q(t)`$. ∎
###### Proposition 3.3.6.
The Carlitz motive $`C`$ is rigid analytically trivial.
###### Proof.
We see that $`C^{}=𝕃`$ with $`𝝈f=(t\theta )f^{(1)}`$ for $`f𝕃`$. Therefore, by Lemma 3.3.5,
$$C^\mathrm{B}=\{f𝕃(t\theta )f^{(1)}=f\}=\frac{1}{\mathrm{\Omega }}𝔽_q(t).$$
Therefore $`𝕃_{𝔽_q(t)}C^\mathrm{B}C^{}`$. ∎
###### Lemma 3.3.7.
Let $`P`$ be a pre-$`t`$-motive, and let $`\mu _1,\mathrm{},\mu _mP^\mathrm{B}`$. If $`\mu _1,\mathrm{},\mu _m`$ are linearly independent over $`𝔽_q(t)`$, then they are linearly independent over $`𝕃`$ in $`P^{}`$.
###### Proof.
Suppose that $`m2`$ is minimal such that $`\mu _1,\mathrm{},\mu _m`$ are linearly independent over $`𝔽_q(t)`$ but that $`_{i=1}^mf_i\mu _i=0`$, with $`f_i𝕃,f_1=1`$. Now,
$$𝝈\underset{i=1}{\overset{m}{}}f_i\mu _i=\underset{i=1}{\overset{m}{}}f_i^{(1)}\mu _i=0.$$
Therefore, $`_{i=2}^m(f_if_i^{(1)})\mu _i=0`$. By the minimality of $`m`$ and Lemma 3.3.2, each $`f_i`$ is in $`𝔽_q(t)`$. However, this violates the $`𝔽_q(t)`$-linear independence of $`\mu _1,\mathrm{},\mu _m`$. ∎
###### Proposition 3.3.8.
If $`P`$ is a pre-$`t`$-motive, then $`dim_{𝔽_q(t)}P^\mathrm{B}dim_{\overline{k}(t)}P`$. Equality holds if and only if $`P`$ is rigid analytically trivial.
###### Proof.
From Lemma 3.3.7, the map $`𝕃_{𝔽_q(t)}P^\mathrm{B}P^{}`$ is injective. The inequality in the statement of the proposition follows from the equality $`dim_{\overline{k}(t)}P=dim_𝕃P^{}`$. By the definition of rigid analytic triviality, equality holds if and only if the map above is also surjective. ∎
###### Proposition 3.3.9.
Suppose that $`P`$ is a pre-$`t`$-motive and that $`\mathrm{\Phi }`$ represents multiplication by $`𝛔`$ on $`P`$ with respect to the basis $`𝐩`$ of $`P`$.
1. $`P`$ is rigid analytically trivial if and only if there is a matrix $`\mathrm{\Psi }\mathrm{GL}_r(𝕃)`$ satisfying
$$\sigma (\mathrm{\Psi })=\mathrm{\Psi }^{(1)}=\mathrm{\Phi }\mathrm{\Psi }.$$
Such a matrix $`\mathrm{\Psi }`$ is called a *rigid analytic trivialization of $`\mathrm{\Phi }`$* (cf. \[1, Thm. 5\], \[2, Lem. 4.4.13\]).
2. If $`\mathrm{\Psi }`$ is a rigid analytic trivialization of $`\mathrm{\Phi }`$, then the entries of $`\mathrm{\Psi }^1𝐩`$ form an $`𝔽_q(t)`$-basis for $`P^\mathrm{B}`$.
3. If $`P`$ is rigid analytically trivial, $`\mathrm{\Phi }\mathrm{Mat}_r(\overline{k}[t])`$, and $`det(\mathrm{\Phi })=d(t\theta )^s`$ for some $`s0`$ and $`d\overline{k}^\times `$, then there is a rigid analytic trivialization $`\mathrm{\Psi }`$ of $`\mathrm{\Phi }`$ with $`\mathrm{\Psi }\mathrm{GL}_r(𝕋)`$.
###### Proof.
The proofs of parts (a) and (b) are essentially the same as the proof of \[2, Lem. 4.4.13\] with minor modifications. We provide a sketch for completeness. ((a) $``$; (b)): Certainly if we have such a $`\mathrm{\Psi }`$, then the entries of $`\mathrm{\Psi }^1𝐩`$ are both an $`𝕃`$-basis of $`P^{}`$ and also an $`𝔽_q(t)`$-linearly independent set in $`P^\mathrm{B}`$. By Proposition 3.3.8, the entries of $`\mathrm{\Psi }^1𝐩`$ must be an $`𝔽_q(t)`$-basis of $`P^\mathrm{B}`$, and thus $`P`$ is rigid analytically trivial. ((a) $``$): On the other hand, if $`P`$ is rigid analytically trivial, then there is a matrix $`\mathrm{\Theta }\mathrm{GL}_r(𝕃)`$ so that the entries of $`\mathrm{\Theta }𝐩`$ are both an $`𝕃`$-basis of $`P^{}`$ and an $`𝔽_q(t)`$-basis of $`P^\mathrm{B}`$. Setting $`\mathrm{\Psi }:=\mathrm{\Theta }^1`$ gives the desired matrix.
For part (c), we first let $`𝖯`$ be the $`\overline{k}[t]`$-span of the entries of $`𝐩`$, and set
$$𝖯^{}:=𝕋_{\overline{k}[t]}𝖯,𝖯^\mathrm{B}:=\{\mu 𝖯^{}𝝈\mu =\mu \}.$$
For $`\mu P^\mathrm{B}`$, write $`\mu =f_ip_i=𝐟𝐩`$ with $`𝐟\mathrm{Mat}_{1\times r}(𝕃)`$. We claim that for some $`c𝔽_q[t]`$, we have $`c\mu 𝖯^\mathrm{B}`$. Let $`\mathrm{den}(𝐟)𝕂[t]`$ denote the monic least common multiple of the denominators of $`𝐟`$, which is well-defined by (2.2.4.1). Then since $`𝝈\mu =\mu `$, we have
$$𝐟𝐩=𝝈(𝐟𝐩)=𝐟^{(1)}\mathrm{\Phi }𝐩.$$
Therefore, $`\mathrm{den}(𝐟)=\mathrm{den}(𝐟^{(1)}\mathrm{\Phi })`$. But $`\mathrm{\Phi }\mathrm{Mat}_r(\overline{k}[t])`$, so $`\mathrm{den}(𝐟^{(1)}\mathrm{\Phi })`$ divides $`\mathrm{den}(𝐟^{(1)})`$. Degree considerations force $`\mathrm{den}(𝐟)=\mathrm{den}(𝐟^{(1)})`$. Therefore take $`c=\mathrm{den}(𝐟)𝔽_q[t]`$. This proves the claim, and moreover we have shown that
$$P^\mathrm{B}𝔽_q(t)_{𝔽_q[t]}𝖯^\mathrm{B}.$$
Furthermore, it follows that as $`𝕃`$-vector spaces, $`P^{}𝕃_𝕋𝖯^{}𝕃_{𝔽_q[t]}𝖯^\mathrm{B}`$. Let $`𝝂=[\nu _1,\mathrm{},\nu _r]^{\text{tr}}`$ be an $`𝔽_q[t]`$-basis for $`𝖯^\mathrm{B}`$. Then for some $`\mathrm{\Theta }\mathrm{GL}_r(𝕃)\mathrm{Mat}_r(𝕋)`$, we have $`𝝂=\mathrm{\Theta }𝐩`$. Since $`𝝈𝝂=𝝂`$, it follows that $`\mathrm{\Theta }^{(1)}\mathrm{\Phi }=\mathrm{\Theta }`$. By our initial hypotheses, $`d(t\theta )^sdet(\mathrm{\Theta })^{(1)}=det(\mathrm{\Theta })`$. Choose $`b\overline{k}^\times `$ so that $`d=b^{(1)}/b`$. Then from Lemmas 3.3.2 and 3.3.5 (and the fact that $`\mathrm{\Theta }\mathrm{Mat}_r(𝕋)`$), we see that
$$bdet(\mathrm{\Theta })=\frac{\gamma }{\mathrm{\Omega }^s},\gamma 𝔽_q[t].$$
We claim that $`\gamma 𝔽_q^\times `$. If not, then $`det(\mathrm{\Theta })0(mod\gamma )`$ in $`𝕋`$, and so there is a $`𝐟=[f_1,\mathrm{},f_r]\mathrm{Mat}_{1\times r}(𝕋)`$ so that
$$𝐟\mathrm{\Theta }0(mod\gamma ).$$
Since $`𝕋/\gamma 𝕋𝕂[t]/\gamma 𝕂[t]`$, without loss of generality we can assume that each $`f_i`$ is a polynomial in $`𝕂[t]`$ of degree strictly less than the degree of $`\gamma `$, that $`f_i1`$ for all $`i`$, and that at least one $`f_i`$ satisfies $`f_i=1`$. Now define a norm $`_{}`$ on $`𝖯^{}`$ by
$$h_ip_i_{}:=suph_i,h_1,\mathrm{},h_r𝕋.$$
Then $`_{}`$ defines a complete ultrametic norm on $`𝖯^{}`$ that satisfies
$$h\mu _{}=h\mu _{},h𝕋,\mu 𝖯^{}.$$
Consider
$$𝐟\mathrm{\Theta }^{(1)}\mathrm{\Phi }=𝐟\mathrm{\Theta }0(mod\gamma ).$$
Since $`\gamma `$ is relatively prime to $`det(\mathrm{\Phi })`$, it follows that $`\mathrm{\Phi }`$ is invertible modulo $`\gamma `$, and so
$$𝐟\mathrm{\Theta }^{(1)}0(mod\gamma ).$$
Repeating this argument we find that $`𝐟\mathrm{\Theta }^{(n)}0(mod\gamma )`$ for all $`n0`$. Now, by choice of $`𝐟`$,
$$\frac{1}{\gamma }𝐟𝝂=\frac{1}{\gamma }𝐟\mathrm{\Theta }𝐩𝖯^{},$$
and for each $`n`$, the above congruences for $`𝐟\mathrm{\Theta }^{(n)}`$ imply that $`\gamma ^1𝐟^{(n)}𝝂=\gamma ^1𝐟^{(n)}\mathrm{\Theta }𝐩𝖯^{}`$. Now by Lemma 2.2.7, there is an $`m>0`$ so that with respect to the $`_{}`$ metric,
$$\underset{n\mathrm{}}{lim}\frac{1}{\gamma }f_i^{(mn)}\nu _i=\frac{1}{\gamma }c_i\nu _i𝖯^{},$$
where $`c_i\overline{𝔽}_q[t]`$ and at least one $`c_i0`$, say $`c_a0`$. Now for some $`l1`$, we have every $`c_i𝔽_{q^l}[t]`$. Since the trace map $`𝔽_{q^l}𝔽_q`$ is not trivial, by dividing each $`c_i`$ by a fixed element in $`𝔽_{q^l}^\times `$, we can assume that $`c_a+c_a^{(1)}+\mathrm{}+c_a^{(1l)}0`$. Therefore,
$$\underset{j=0}{\overset{l1}{}}𝝈^j\left(\frac{1}{\gamma }\underset{i=1}{\overset{r}{}}c_i\nu _i\right)=\frac{1}{\gamma }\underset{i=1}{\overset{r}{}}\left(\underset{j=0}{\overset{l1}{}}c_i^{(j)}\right)\nu _i𝖯^{}.$$
Thus we obtain $`\mu :=\gamma ^1d_i\nu _i𝖯^{}`$, $`d_i𝔽_q[t]`$, $`d_a0`$. Easily we see that $`\mu 𝖯^\mathrm{B}`$ and $`\mu 0`$. Since $`\mathrm{deg}f_i<\mathrm{deg}\gamma `$ for each $`i`$, we have $`\mathrm{deg}d_i<\mathrm{deg}\gamma `$ for each $`i`$. In particular, $`\gamma `$ does not divide $`d_a`$. Thus $`\mu 𝖯^\mathrm{B}`$ but $`\mu `$ is not in the $`𝔽_q[t]`$-span of $`𝝂`$, which contradicts that $`𝝂`$ is an $`𝔽_q[t]`$-basis of $`𝖯^\mathrm{B}`$. Therefore, it follows that $`\gamma 𝔽_q^\times `$, and since $`\mathrm{\Omega }𝕋^\times `$, we have $`det(\mathrm{\Theta })𝕋^\times `$. Taking $`\mathrm{\Psi }=\mathrm{\Theta }^1`$ provides the desired rigid analytic trivialization. ∎
#### 3.3.10. Remark
It is worth noting that multiplication by $`\mathrm{\Theta }`$ induces the isomorphism of $`𝕃`$-vector spaces,
$$𝕃_𝕋𝖯^{}𝕃_𝕋(𝕋_{𝔽_q[t]}𝖯^\mathrm{B}).$$
Since $`\mathrm{\Theta }\mathrm{GL}_r(𝕋)`$, this then implies $`𝖯^{}𝕋_{𝔽_q[t]}𝖯^\mathrm{B}`$ as $`𝕋`$-modules.
###### Proposition 3.3.11.
Let
$$0PQR0$$
be an exact sequence of pre-$`t`$-motives.
1. If $`Q`$ is rigid analytically trivial, then both $`P`$ and $`R`$ are rigid analytically trivial.
2. If $`P`$, $`Q`$, and $`R`$ are rigid analytically trivial, then the sequence
$$0P^\mathrm{B}Q^\mathrm{B}R^\mathrm{B}0$$
is an exact sequence of $`𝔽_q(t)`$-vector spaces.
###### Proof.
The sequence $`0P^\mathrm{B}Q^\mathrm{B}R^\mathrm{B}`$ is exact. Now suppose that $`Q`$ is rigid analytically trivial. Let $`\kappa :𝕃_{𝔽_q(t)}Q^\mathrm{B}𝕃_{𝔽_q(t)}R^\mathrm{B}`$ be the natural map. Then we have a commutative diagram with exact rows,
where the central vertical map is an isomorphism by hypothesis, and the other two are injective by Lemma 3.3.7. The injectivity of all three maps then implies that each is an isomorphism. Thus we see immediately that $`P`$ is rigid analytically trivial. Also we see that
$$dim_{𝔽_q(t)}R^\mathrm{B}=dim_𝕃𝕃_{𝔽_q(t)}R^\mathrm{B}dim_𝕃\mathrm{im}(\kappa )=dim_𝕃R^{}=dim_{\overline{k}(t)}R,$$
which by Proposition 3.3.8 must be a string of equalities. Therefore $`R`$ is rigid analytically trivial, which completes part (a).
Now suppose that $`P`$, $`Q`$, and $`R`$ are all rigid analytically trivial. Then
$$dim_{𝔽_q(t)}Q^\mathrm{B}=dim_{\overline{k}(t)}Q=dim_{\overline{k}(t)}P+dim_{\overline{k}(t)}R=dim_{𝔽_q(t)}P^\mathrm{B}+dim_{𝔽_q(t)}R^\mathrm{B},$$
which proves part (b). ∎
#### 3.3.12. Remark
In particular, it follows from Proposition 3.3.11 that kernels and cokernels exist in $``$, which implies that $``$ is an abelian $`𝔽_q(t)`$-linear category. We also see that
$$PP^\mathrm{B}:\mathrm{𝐕𝐞𝐜}(𝔽_q(t)),$$
where $`\mathrm{𝐕𝐞𝐜}(𝔽_q(t))`$ is the category of finite dimensional vector spaces over $`𝔽_q(t)`$, is an exact $`𝔽_q(t)`$-linear functor.
###### Proposition 3.3.13.
Let $`P`$ and $`Q`$ be rigid analytically trivial pre-$`t`$-motives. Then the natural map
$$\mathrm{Hom}_{}(P,Q)\mathrm{Hom}_{𝔽_q(t)}(P^\mathrm{B},Q^\mathrm{B})$$
is injective.
###### Proof.
Suppose $`\varphi :PQ`$ is a morphism in $`\mathrm{Hom}_{}(P,Q)`$. Then we have an exact sequence in $``$,
$$0\mathrm{ker}\varphi P\stackrel{\varphi }{}QQ/\varphi (P)0,$$
which leads then to an exact sequence of $`𝔽_q(t)`$-vector spaces,
$$0(\mathrm{ker}\varphi )^\mathrm{B}P^\mathrm{B}\stackrel{\varphi ^\mathrm{B}}{}Q^\mathrm{B}(Q/\varphi (P))^\mathrm{B}0.$$
Since the dimension over $`\overline{k}(t)`$ of each term in the first sequence is the same as the dimension over $`𝔽_q(t)`$ of the corresponding term in the second sequence, we see that $`\varphi ^\mathrm{B}=0`$ if and only if $`\varphi =0`$. ∎
###### Proposition 3.3.14.
If pre-$`t`$-motives $`P`$ and $`Q`$ are rigid analytically trivial, then
1. $`PQ`$ is rigid analytically trivial, and the natural map $`P^\mathrm{B}_{𝔽_q(t)}Q^\mathrm{B}(PQ)^\mathrm{B}`$ is an isomorphism of $`𝔽_q(t)`$-vector spaces;
2. $`P^{}`$ is rigid analytically trivial, and the natural map $`\left(P^\mathrm{B}\right)^{}\left(P^{}\right)^\mathrm{B}`$ is an isomorphism of $`𝔽_q(t)`$-vector spaces.
###### Proof.
Here we make use of Proposition 3.3.9. We first note that
$$(PQ)^{}=𝕃_{\overline{k}(t)}(P_{\overline{k}(t)}Q)(𝕃_{\overline{k}(t)}P)_𝕃(𝕃_{\overline{k}(t)}Q)=P^{}_𝕃Q^{},$$
where the middle isomorphism is an isomorphism of $`𝕃`$-vector spaces that commutes with the action of $`𝝈`$. We observe that we can choose $`\overline{k}(t)`$-bases for $`P`$, $`Q`$, and $`PQ`$ so that multiplication by $`𝝈`$ is represented by matrices $`\mathrm{\Phi }_P`$, $`\mathrm{\Phi }_Q`$, and $`\mathrm{\Phi }_{PQ}`$ satisfying
$$\mathrm{\Phi }_{PQ}=\mathrm{\Phi }_P\mathrm{\Phi }_Q.$$
By Proposition 3.3.9, we can choose $`\mathrm{\Psi }_P`$, $`\mathrm{\Psi }_Q\mathrm{GL}_r(𝕃)`$ that are rigid analytic trivializations of $`\mathrm{\Phi }_P`$ and $`\mathrm{\Phi }_Q`$. Then we note that $`\mathrm{\Psi }_{PQ}:=\mathrm{\Psi }_P\mathrm{\Psi }_Q`$ is a rigid analytic trivialization of $`\mathrm{\Phi }_{PQ}`$. Now note that $`\mathrm{\Phi }_P^{}:=(\mathrm{\Phi }_P^1)^{\text{tr}}`$ represents multiplication by $`𝝈`$ with respect to the dual basis and that $`\mathrm{\Psi }_P^{}:=(\mathrm{\Psi }_P^1)^{\text{tr}}`$ is a rigid analytic trivialization. The second parts of (a) and (b) are straightforward. ∎
###### Theorem 3.3.15.
The category $``$ of rigid analytically trivial pre-$`t`$-motives is a neutral Tannakian category over $`𝔽_q(t)`$ with fiber functor $`PP^\mathrm{B}:\mathrm{𝐕𝐞𝐜}(𝔽_q(t))`$.
###### Proof.
We have seen that
* $`\mathrm{𝟏}`$ is in $``$ (Proposition 3.3.3);
* $``$ is an abelian category (Proposition 3.3.11 and §3.3.12);
* $``$ is closed under tensor products and duals (Proposition 3.3.14).
Thus $``$ is a rigid abelian $`𝔽_q(t)`$-linear tensor subcategory of $`𝒫`$ (see \[11, Defs. II.1.14-15\]). We have also shown that
* $`\mathrm{End}_{}(\mathrm{𝟏})=𝔽_q(t)`$ (Lemma 3.2.9);
* For each $`P`$ in $``$, the $`𝔽_q(t)`$-vector space $`P^\mathrm{B}`$ is finite dimensional (Proposition 3.3.8);
* $`PP^\mathrm{B}`$ is $`𝔽_q(t)`$-linear and exact (Proposition 3.3.11 and §3.3.12);
* $`PP^\mathrm{B}`$ is faithful (Proposition 3.3.13);
* $`PP^\mathrm{B}`$ is a tensor functor (Proposition 3.3.14).
Thus $``$ is a neutral Tannakian category over $`𝔽_q(t)`$ with fiber functor $`PP^\mathrm{B}`$ (see \[11, Def. II.2.19\]). ∎
### 3.4. Anderson $`t`$-motives
Here we recall the definitions and essential properties of “dual $`t`$-motives” from . So as not to confuse these objects with the duals of $`t`$-motives to be used later on, we call these objects *Anderson $`t`$-motives*, since they are simply the dual notion of the objects studied in .
#### 3.4.1. Definition
An *Anderson $`t`$-motive* $`𝖬`$ is a left $`\overline{k}[t;𝝈]`$-module such that
* $`𝖬`$ is free and finitely generated over $`\overline{k}[t]`$;
* $`𝖬`$ is free and finitely generated over $`\overline{k}[𝝈]`$;
* $`(t\theta )^n𝖬𝝈𝖬`$ for all $`n0`$.
A morphism of Anderson $`t`$-motives is a left $`\overline{k}[t;𝝈]`$-module homomorphism. In this way Anderson $`t`$-motives form a category.
As in §3.2.3, if $`𝗆\mathrm{Mat}_{r\times 1}(𝖬)`$ is a $`\overline{k}[t]`$-module basis for $`𝖬`$, then there is a matrix $`\mathrm{\Phi }=\mathrm{\Phi }_𝗆\mathrm{Mat}_{r\times 1}(\overline{k}[t])`$ so that
$$𝝈𝗆=\mathrm{\Phi }𝗆.$$
Since a power of $`t\theta `$ annihilates $`𝖬/𝝈𝖬`$, we have
$$det\mathrm{\Phi }=c(t\theta )^s$$
for some $`c\overline{k}^\times `$, where $`s`$ is the rank of $`𝖬`$ as a $`\overline{k}[𝝈]`$-module.
#### 3.4.2. Anderson $`t`$-motives to pre-$`t`$-motives
Given an Anderson $`t`$-motive $`𝖬`$ we obtain a pre-$`t`$-motive $`M`$ by setting
$$M:=\overline{k}(t)_{\overline{k}[t]}𝖬$$
and defining
$$𝝈(fm):=f^{(1)}𝝈m.$$
It is straightforward to check that $`M`$ is a left $`\overline{k}(t)[𝝈,𝝈^1]`$-module, and it is of course finite dimensional as a $`\overline{k}(t)`$-vector space. Moreover, $`𝖬M`$ is a functor from the category of Anderson $`t`$-motives to the category of pre-$`t`$-motives.
#### 3.4.3. The Carlitz motive
Let $`𝖢`$ be the Anderson $`t`$-motive whose underlying $`\overline{k}[t]`$-module is $`\overline{k}[t]`$ itself. Then the action of $`𝝈`$ on $`𝖢`$ is defined by
$$𝝈(f)=(t\theta )f^{(1)},f𝖢.$$
It is not difficult to check that $`𝖢`$ is an Anderson $`t`$-motive, and that its image in $`𝒫`$ is the Carlitz motive. For any $`n1`$, we also have the $`n`$-th tensor power of $`𝖢`$,
$$𝖢(n):=𝖢_{\overline{k}[t]}\mathrm{}_{\overline{k}[t]}𝖢,$$
with diagonal $`𝝈`$-action. It is an Anderson $`t`$-motive sent to $`C(n)`$ in $`𝒫`$.
#### 3.4.4. The Carlitz module
The *Carlitz module $``$ over $`\overline{k}`$* is defined to be the $`𝔽_q`$-algebra $`\overline{k}`$ together with an $`𝔽_q[t]`$-module structure defined by
$$_t(x):=\theta x+x^q,x\overline{k}.$$
That is the $`𝔽_q`$-algebra homomorphism $`a_a:𝔽_q[t]\overline{k}[𝝈^1]`$ defined by $`t\theta +𝝈^1`$ induces an $`𝔽_q[t]`$-module structure on $`\overline{k}`$. See \[15, Ch. 3\] or \[31, §2.5\] for more details. To see the relationship with the Carlitz motive, we note that there is an isomorphism
$$(\overline{k})\frac{𝖢}{(𝝈1)𝖢}$$
of $`𝔽_q[t]`$-modules. Indeed if $`x\overline{k}`$, then
$$tx=\theta x+(t\theta )x=\theta x+𝝈(x^q)=\theta x+x^q+(𝝈1)x^q.$$
Similarly $`ax_a(x)(mod𝝈1)`$ for all $`a𝔽_q[t]`$. It is a simple matter to check that there is a natural isomorphism of $`𝔽_q`$-vector spaces $`𝖢/(𝝈1)𝖢\overline{k}`$. Thus $`𝖢/(𝝈1)𝖢`$ presents the Carlitz module directly.
###### Proposition 3.4.5.
For Anderson $`t`$-motives $`𝖬`$ and $`𝖭`$, the natural map
$$\mathrm{Hom}_{\overline{k}[t;𝝈]}(𝖬,𝖭)_{𝔽_q[t]}𝔽_q(t)\mathrm{Hom}_𝒫(M,N)$$
is an isomorphism of $`𝔽_q(t)`$-vector spaces.
###### Proof.
Let $`\mathrm{\Theta }`$ denote the map in question. It is clearly $`𝔽_q(t)`$-linear. To see that it is injective, we first observe that if $`\alpha \mathrm{Hom}_{\overline{k}[t;𝝈]}(𝖬,𝖭)_{𝔽_q[t]}𝔽_q(t)`$ then $`\alpha =\underset{¯}{\varphi }\frac{1}{v}`$, for some $`\underset{¯}{\varphi }\mathrm{Hom}_{\overline{k}[t;𝝈]}(𝖬,𝖭)`$ and $`v𝔽_q[t]`$, $`v0`$. Then $`v\mathrm{\Theta }(\alpha )=\mathrm{\Theta }(v\alpha )=\mathrm{\Theta }(\underset{¯}{\varphi }1)=:\varphi `$. But
$`\underset{¯}{\varphi }`$ $`\mathrm{Hom}_{\overline{k}[t;𝝈]}(𝖬,𝖭)\mathrm{Hom}_{\overline{k}[t]}(𝖬,𝖭),`$
$`\varphi `$ $`\mathrm{Hom}_{\overline{k}(t)[𝝈,𝝈^1]}(M,N)\mathrm{Hom}_{\overline{k}(t)}(M,N),`$
and so $`\varphi =0`$ if and only if $`\underset{¯}{\varphi }=0`$. Thus $`\mathrm{\Theta }(\alpha )=0`$ if and only if $`\alpha =0`$.
For surjectivity, suppose that $`\varphi \mathrm{Hom}_𝒫(M,N)`$. Fix $`\overline{k}[t]`$-bases $`𝗆`$ and $`𝗇`$ for $`𝖬`$ and $`𝖭`$ respectively, and extend these to bases $`𝐦`$ and $`𝐧`$ of $`M`$ and $`N`$. Then the map $`\varphi :MN`$ is represented by a matrix $`F\mathrm{Mat}_{r\times s}(\overline{k}(t))`$ so that $`F^{(1)}\mathrm{\Phi }_𝐧=\mathrm{\Phi }_𝐦F`$ as in §3.2.3. By choice of $`𝗆`$ and $`𝗇`$, $`\mathrm{\Phi }_𝐦`$ and $`\mathrm{\Phi }_𝐧`$ have entries in $`\overline{k}[t]`$, and it suffices to show that $`F`$ has entries with denominators in $`𝔽_q[t]`$.
For a matrix $`B`$ with entries in $`\overline{k}(t)`$, let $`\mathrm{den}(B)\overline{k}[t]`$ be the monic least common multiple of the denominators of the entries of $`B`$. Since $`det(\mathrm{\Phi }_𝐧)=c(t\theta )^s`$ for some $`s0`$ and $`c\overline{k}^\times `$, we see that
$$\mathrm{den}(F)(t\theta )^sF^{(1)}=\mathrm{den}(F)(t\theta )^s\mathrm{\Phi }_𝐦F\mathrm{\Phi }_𝐧^1\mathrm{Mat}_{r\times s}(\overline{k}[t]).$$
Therefore, $`\mathrm{den}(F^{(1)})`$ divides $`\mathrm{den}(F)(t\theta )^s`$. However, $`\mathrm{den}(F^{(1)})=\mathrm{den}(F)^{(1)}`$ and so $`\mathrm{deg}(\mathrm{den}(F^{(1)}))=\mathrm{deg}(\mathrm{den}(F))`$. Thus, it suffices to show that $`\mathrm{den}(F^{(1)})`$ is relatively prime to $`t\theta `$, since then $`\mathrm{den}(F)^{(1)}=\mathrm{den}(F)`$ whence all of the denominators of $`F`$ are in $`𝔽_q[t]`$.
Suppose that $`t\theta `$ divides $`\mathrm{den}(F^{(1)})`$, and so $`t\theta ^q`$ divides $`\mathrm{den}(F)`$. Then $`t\theta ^q`$ divides $`\mathrm{den}(\mathrm{\Phi }_𝐦F)`$, because otherwise $`t\theta ^q`$ would divide $`det(\mathrm{\Phi }_𝐦)`$ which is a power of $`t\theta `$. Likewise, $`t\theta ^q`$ divides $`\mathrm{den}(\mathrm{\Phi }_𝐦F\mathrm{\Phi }_𝐧^1)=\mathrm{den}(F^{(1)})`$. By repeating the same argument we see that $`\mathrm{den}(F^{(1)})`$ is divisible by each of
$$t\theta ,t\theta ^q,t\theta ^{q^2},\mathrm{}$$
contradicting that $`\mathrm{den}(F^{(1)})\overline{k}[t]`$. ∎
#### 3.4.6. Rigid analytic triviality
Similar to §3.3.1, if $`𝖬`$ is an Anderson $`t`$-motive, then we set
$$𝖬^{}:=𝕋_{\overline{k}[t]}𝖬.$$
We provide $`𝖬^{}`$ with a $`\overline{k}[t;𝝈]`$-module structure by setting $`𝝈(fm)=f^{(1)}𝝈m`$, and we set
$$𝖬^\mathrm{B}:=(𝖬^{})^𝝈=\{\mu 𝖬^{}𝝈\mu =\mu \}.$$
We say that $`𝖬`$ is *rigid analytically trivial* if the natural map $`𝕋_{𝔽_q[t]}𝖬^\mathrm{B}𝖬^{}`$ is an isomorphism. The following proposition is a companion to Proposition 3.3.9.
###### Proposition 3.4.7.
Let $`𝖬`$ be an Anderson $`t`$-motive, and let $`M`$ be its corresponding pre-$`t`$-motive. Suppose $`𝗆\mathrm{Mat}_{r\times 1}(𝖬)`$ is a $`\overline{k}[t]`$-basis for $`𝖬`$, and let $`\mathrm{\Phi }\mathrm{Mat}_{r\times 1}(\overline{k}[t])`$ represent multiplication by $`𝛔`$ on $`𝖬`$ with respect to $`𝗆`$.
1. $`𝖬`$ is rigid analytically trivial if and only if it admits a rigid analytic trivialization $`\mathrm{\Psi }`$ with $`\mathrm{\Psi }\mathrm{GL}_r(𝕋)`$.
2. If $`\mathrm{\Psi }\mathrm{GL}_r(𝕋)`$ is a rigid analytic trivialization of $`\mathrm{\Phi }`$, then the entries of $`\mathrm{\Psi }^1𝗆`$ form an $`𝔽_q[t]`$-basis of $`𝖬^\mathrm{B}`$.
3. $`𝖬`$ is rigid analytically trivial if and only if $`M`$ is rigid analytically trivial.
###### Proof.
The proofs of parts (a) and (b) are in \[2, Lem. 4.4.13\] and follow the same lines as their counterparts in Proposition 3.3.9. Part (c) is then a consequence of Proposition 3.3.9(c). ∎
#### 3.4.8. Definition
We define the *category $`𝒜^I`$ of Anderson $`t`$-motives up to isogeny* as follows:
* Objects of $`𝒜^I`$: Anderson $`t`$-motives;
* Morphisms of $`𝒜^I`$: For Anderson $`t`$-motives $`𝖬`$ and $`𝖭`$,
$$\mathrm{Hom}_{𝒜^I}(𝖬,𝖭):=\mathrm{Hom}_{\overline{k}[t;𝝈]}(𝖬,𝖭)_{𝔽_q[t]}𝔽_q(t).$$
We also define the full subcategory $`𝒜^I`$ of rigid analytically trivial Anderson $`t`$-motives up to isogeny by restriction. We sum up the results of this section in the following theorem.
###### Theorem 3.4.9.
Let $`𝒫`$ be the category of pre-$`t`$-motives, and let $``$ be the category of rigid analytically trivial pre-$`t`$-motives.
1. The functor $`𝖬M:𝒜^I𝒫`$ is fully faithful.
2. The functor $`𝖬M:𝒜^I`$ is fully faithful.
###### Proof.
Part (a) is simply a restatement of Proposition 3.4.5. That the functor in part (b) is well-defined follows from Proposition 3.4.7(c), and its full faithfulness follows from Proposition 3.4.5. ∎
#### 3.4.10. The category $`𝒯`$
We define the *category $`𝒯`$ of $`t`$-motives* to be the strictly full Tannakian subcategory of $``$ generated by the essential image of the functor
$$𝖬M:𝒜^I.$$
The category of $`t`$-motives can further be described as follows:
* Objects of $`𝒯`$: rigid analytically trivial pre-$`t`$-motives that can be constructed from Anderson $`t`$-motives using direct sums, subquotients, tensor products, duals, and internal Hom’s.
* Morphisms of $`𝒯`$: morphisms of left $`\overline{k}(t)[𝝈,𝝈^1]`$-modules.
It is worth noting that Proposition 3.4.7(c) says that the category of $`t`$-motives is the strictly full Tannakian subcategory of $``$ generated by the intersection in $`𝒫`$ of $``$ and the image of *all* Anderson $`t`$-motives.
### 3.5. Galois groups of $`t`$-motives
Having defined a Tannakian category of $`t`$-motives, it is now possible to assign to each $`t`$-motive a linear algebraic group over $`𝔽_q(t)`$, which we call the *Galois group* of the $`t`$-motive. For essential facts about Tannakian categories and their associated groups, we refer to , , \[26, App. B\].
#### 3.5.1. Fiber functors
The functor
$$\begin{array}{ccc}\hfill \omega :𝒯& & \mathrm{𝐕𝐞𝐜}(𝔽_q(t))\hfill \\ \hfill M& & M^\mathrm{B}\hfill \end{array}$$
is the fiber functor of $`𝒯`$. For any commutative $`𝔽_q(t)`$-algebra $`R`$, we let $`\omega ^{(R)}:𝒯\mathrm{𝐌𝐨𝐝}(R)`$ be the extension of $`\omega `$ defined by
$$\omega ^{(R)}(M):=R_{𝔽_q(t)}M^\mathrm{B},$$
where $`\mathrm{𝐌𝐨𝐝}(R)`$ is the category of finitely generated left $`R`$-modules. Now fix a $`t`$-motive $`M`$. We let $`𝒯_M`$ be the strictly full Tannakian subcategory of $`𝒯`$ generated by $`M`$. That is, $`𝒯_M`$ consists of all objects of $`𝒯`$ isomorphic to subquotients of finite direct sums of $`M^u(M^{})^v`$ for various $`u`$, $`v`$. The fiber functor of $`𝒯_M`$ is $`\omega _M:𝒯_M\mathrm{𝐕𝐞𝐜}(𝔽_q(t))`$, the restriction of $`\omega `$ to $`𝒯_M`$, and similarly we restrict $`\omega _M^{(R)}`$ to $`𝒯_M`$ for an $`𝔽_q(t)`$-algebra $`R`$.
#### 3.5.2. Galois groups
As $`𝒯`$ is a neutral Tannakian category over $`𝔽_q(t)`$, there is an affine group scheme $`\mathrm{\Gamma }_𝒯`$ over $`𝔽_q(t)`$ so that $`𝒯`$ is equivalent to the category $`\mathrm{𝐑𝐞𝐩}(\mathrm{\Gamma }_𝒯,𝔽_q(t))`$ of finite dimensional representations of $`\mathrm{\Gamma }_𝒯`$ over $`𝔽_q(t)`$:
$$𝒯\mathrm{𝐑𝐞𝐩}(\mathrm{\Gamma }_𝒯,𝔽_q(t)).$$
The group $`\mathrm{\Gamma }_𝒯`$ is defined to be the group of tensor automorphisms of the fiber functor $`\omega `$; that is, if $`R`$ is any $`𝔽_q(t)`$-algebra, then
$$\mathrm{\Gamma }_𝒯(R)=\mathrm{Aut}_𝒯^{}\left(\omega ^{(R)}\right).$$
Now for any $`t`$-motive $`M`$, there is a linear algebraic group $`\mathrm{\Gamma }_M:=\mathrm{\Gamma }_{𝒯_M}`$ over $`𝔽_q(t)`$ so that $`𝒯_M`$ is equivalent to $`\mathrm{𝐑𝐞𝐩}(\mathrm{\Gamma }_M,𝔽_q(t))`$. As such, for any $`𝔽_q(t)`$-algebra $`R`$, $`\mathrm{\Gamma }_M(R)=\mathrm{Aut}_{𝒯_M}^{}\left(\omega _M^{(R)}\right)`$. In this way we find that we have a naturally defined faithful representation
$$\mathrm{\Gamma }_M\mathrm{GL}(M^\mathrm{B})$$
over $`𝔽_q(t)`$, which provides the basis for constructing the equivalence of categories,
$$𝒯_M\mathrm{𝐑𝐞𝐩}(\mathrm{\Gamma }_M,𝔽_q(t)).$$
The group $`\mathrm{\Gamma }_M`$ is called the *Galois group of $`M`$*. Furthermore, there is a surjective group homomorphism,
$$\mathrm{\Gamma }_𝒯\mathrm{\Gamma }_M.$$
If $`N`$ is another $`t`$-motive in $`𝒯_M`$, then there is a natural surjective homomorphism, $`\mathrm{\Gamma }_M\mathrm{\Gamma }_N`$. In §4, we will show that $`\mathrm{\Gamma }_M`$ can be calculated using systems of $`\sigma `$-semilinear equations. For now we will calculate the Galois group of the Carlitz motive $`C`$.
###### Lemma 3.5.3.
For $`m`$, $`n`$,
$$\mathrm{Hom}_𝒯(C(m),C(n))\{\begin{array}{cc}𝔽_q(t)\hfill & \text{if }m=n\text{,}\hfill \\ 0\hfill & \text{if }mn\text{.}\hfill \end{array}$$
###### Proof.
By tensoring with $`C(m)`$, we see that $`\mathrm{Hom}_𝒯(C(m),C(n))\mathrm{Hom}_𝒯(\mathrm{𝟏},C(nm))`$. Thus it suffices to assume that $`m=0`$. If $`\varphi :\mathrm{𝟏}C(n)`$ is a morphism in $`𝒯`$, then $`\varphi `$ is represented by some $`a\overline{k}(t)^\times `$ such that $`a=a^{(1)}(t\theta )^n`$. By Lemma 3.3.5, this equation has no non-zero solutions $`a\overline{k}(t)`$ unless $`n=0`$, in which case $`a`$ can be anything in $`𝔽_q(t)`$. ∎
###### Theorem 3.5.4.
For the Carlitz motive $`C`$, there is an isomorphism $`\mathrm{\Gamma }_C𝔾_\mathrm{m}`$ over $`𝔽_q(t)`$.
###### Proof.
It is easy enough to check this theorem directly. However, by Lemma 3.5.3, $`𝒯_C`$ is equivalent to a $``$-graded category of vector spaces over $`\overline{k}(t)`$ with a fiber functor to $`\mathrm{𝐕𝐞𝐜}(𝔽_q(t))`$, and so its Galois group is $`𝔾_\mathrm{m}`$ over $`𝔽_q(t)`$ \[11, Ex. II.2.30\]. ∎
## 4. Galois theory of systems of $`\sigma `$-semilinear equations
In this section we demonstrate how to calculate the Galois group of a $`t`$-motive as the Galois group of a system of difference equations with respect to the automorphism $`\sigma :\overline{k}(t)\overline{k}(t)`$. These systems of equations and their Galois groups are similar to systems of linear differential equations and their Galois groups, and one should compare our constructions with \[25, Ch. 1\], \[26, Chs. 1–2\], which we have used as guides, as well as , , , , , . For an example of a Galois group of this type in the context of $`t`$-motives, see also the proof of \[7, Prop. 7.1\].
Van der Put and Singer have developed the theory of Picard-Vessiot rings for linear difference equations which is quite useful in our context. However, their treatment generally assumes that the field of constants is algebraically closed. In our case the field of constants is $`𝔽_q(t)`$, which presents several difficulties. On the other hand, the Picard-Vessiot rings treated in are not always domains, whereas our central Picard-Vessiot rings *are* domains by construction, which provides several benefits for the characterization of their Galois groups. It is worth noting that some of what is covered here is covered by the theory of Y. André , but we present everything from scratch for completeness. We thank the referee for making several useful suggestions for improving the clarity of this section.
### 4.1. Solutions of $`\sigma `$-semilinear equations
#### 4.1.1. Fields of definition
Let $`FKL`$ be fields together with an automorphism $`\sigma :LL`$. We say that the triple $`(F,K,L)`$ is *$`\sigma `$-admissible* if
* $`\sigma `$ restricts to automorphisms of $`F`$ and $`K`$;
* $`F=F^\sigma =K^\sigma =L^\sigma `$;
* $`L`$ is a separable extension of $`K`$.
The primary example of $`\sigma `$-admissible fields that we have in mind is
$$(F,K,L)=(𝔽_q(t),\overline{k}(t),𝕃),$$
with automorphism $`\sigma `$ defined as in §2.2.5 by $`\sigma (f)=f^{(1)}`$. This example will be important for applications to $`t`$-motives in §4.5. To see that this triple is $`\sigma `$-admissible, we know that $`𝔽_q(t)=𝔽_q(t)^\sigma =\overline{k}(t)^\sigma `$ by definition and that $`𝕃^\sigma =𝔽_q(t)`$ by Lemma 3.3.2. Also, since $`𝕃`$ is linearly disjoint from $`\overline{k}(t^{1/p})`$, it is therefore separable over $`\overline{k}(t)`$ \[22, Thm. 26.3\]. Henceforth we shall assume that a $`\sigma `$-admissible triple $`(F,K,L)`$ has been chosen.
#### 4.1.2. Convention
If $`\rho :SR`$ is a homomorphism of modules or rings, and $`B\mathrm{Mat}_{r\times s}(S)`$, we let $`\rho (B)\mathrm{Mat}_{r\times s}(R)`$ be the matrix obtained by applying $`\rho `$ to the entries of $`B`$.
#### 4.1.3. Definition
Given a matrix $`\mathrm{\Phi }\mathrm{GL}_r(K)`$, we consider vectors $`\psi \mathrm{Mat}_{r\times 1}(L)`$ that satisfy
$$\sigma (\psi )=\mathrm{\Phi }\psi .$$
In this way, we define a *system of $`\sigma `$-semilinear equations*, and $`\psi `$ is a solution. The set of solutions
$$\mathrm{Sol}(\mathrm{\Phi }):=\{\psi \mathrm{Mat}_{r\times 1}(L)\sigma (\psi )=\mathrm{\Phi }\psi \}$$
is an $`F`$-vector space.
###### Lemma 4.1.4.
Let $`\mathrm{\Phi }\mathrm{GL}_r(K)`$. Suppose that $`\psi _1,\mathrm{},\psi _m\mathrm{Sol}(\mathrm{\Phi })`$ are linearly independent over $`F`$. Then they are linearly independent over $`L`$.
###### Proof.
The proof is in the same spirit as the one for Lemma 3.3.7, and we omit it. ∎
###### Corollary 4.1.5.
Let $`\mathrm{\Phi }\mathrm{Mat}_r(K)`$. Then $`\mathrm{Sol}(\mathrm{\Phi })`$ is an $`F`$-vector space of dimension at most $`r`$.
#### 4.1.6. Fundamental matrix of solutions
Given $`\mathrm{\Phi }\mathrm{GL}_r(K)`$, suppose $`\mathrm{\Psi }\mathrm{GL}_r(L)`$ satisfies
$$\sigma (\mathrm{\Psi })=\mathrm{\Phi }\mathrm{\Psi }.$$
Then by Lemma 4.1.4 and Corollary 4.1.5, the columns of $`\mathrm{\Psi }`$ form an $`F`$-basis for $`\mathrm{Sol}(\mathrm{\Phi })`$. The matrix $`\mathrm{\Psi }`$ is called a *fundamental matrix for $`\mathrm{\Phi }`$*. It is useful to note that $`\mathrm{\Psi }^{}\mathrm{GL}_r(L)`$ is another fundamental matrix for $`\mathrm{\Phi }`$ if and only if $`\mathrm{\Psi }^1\mathrm{\Psi }^{}`$ is fixed by $`\sigma `$. That is, if and only if $`\mathrm{\Psi }^{}=\mathrm{\Psi }\delta `$ for some $`\delta \mathrm{GL}_r(F)`$.
### 4.2. The difference Galois group
Throughout this section we fix $`\mathrm{\Phi }\mathrm{GL}_r(K)`$ and suppose that $`\mathrm{\Psi }\mathrm{GL}_r(L)`$ is a fundamental matrix for $`\mathrm{\Phi }`$ with respect to our $`\sigma `$-admissible triple $`(F,K,L)`$.
For a ring $`R`$, we let $`\mathrm{GL}_{r/R}`$ denote the $`R`$-group scheme of $`r\times r`$ invertible matrices. Its coordinate ring is $`R[X,1/detX]`$, where $`X=(X_{ij})`$ is an $`r\times r`$ matrix of independent variables. If $`S`$ is an $`R`$-algebra, we will as usual let $`GL_r(S)`$ denote the group of $`S`$-rational points on $`\mathrm{GL}_{r/R}`$. For any $`R`$-scheme $`Z`$, we let $`Z_S:=S\times _RZ`$ be its base extension to an $`S`$-scheme.
#### 4.2.1. Construction of $`\mathrm{\Gamma }`$
We define a $`K`$-algebra map $`\nu :K[X,1/detX]L`$ by setting $`\nu (X_{ij}):=\mathrm{\Psi }_{ij}`$. We let
$$𝔭:=\mathrm{ker}\nu ,\mathrm{\Sigma }:=\mathrm{im}\nu =K[\mathrm{\Psi },1/det\mathrm{\Psi }]L.$$
We let $`\mathrm{\Lambda }`$ be the fraction field of $`\mathrm{\Sigma }`$. Finally, we let $`Z=\mathrm{Spec}\mathrm{\Sigma }`$. In this way $`Z`$ is the small closed subscheme of $`\mathrm{GL}_{r/K}`$ such that $`\mathrm{\Psi }Z(L)`$.
Now set $`\mathrm{\Psi }_1,\mathrm{\Psi }_2\mathrm{GL}_r(L_KL)`$ to be the matrices such that $`(\mathrm{\Psi }_1)_{ij}=\mathrm{\Psi }_{ij}1`$ and $`(\mathrm{\Psi }_2)_{ij}=1\mathrm{\Psi }_{ij}`$, and let $`\stackrel{~}{\mathrm{\Psi }}:=\mathrm{\Psi }_1^1\mathrm{\Psi }_2\mathrm{GL}_r(L_KL)`$. We define an $`F`$-algebra map $`\mu :F[X,1/detX]L_KL`$ by $`\mu (X_{ij})=\stackrel{~}{\mathrm{\Psi }}_{ij}`$. We let
$$𝔮:=\mathrm{ker}\mu ,\mathrm{\Delta }:=\mathrm{im}\mu ,$$
and finally we set $`\mathrm{\Gamma }=\mathrm{Spec}\mathrm{\Delta }`$. In this way $`\mathrm{\Gamma }`$ is the smallest closed subscheme of $`\mathrm{GL}_{r/F}`$ such that $`\stackrel{~}{\mathrm{\Psi }}\mathrm{\Gamma }(L_KL)`$. If we wish denote the dependence on $`\mathrm{\Psi }`$, we will write $`Z_\mathrm{\Psi }`$ and $`\mathrm{\Gamma }_\mathrm{\Psi }`$ for these spaces.
Among other things, we will see in Theorem 4.2.11 that $`\mathrm{\Gamma }`$ is a closed subgroup of $`\mathrm{GL}_{r/F}`$ and that $`Z`$ is a $`\mathrm{\Gamma }_K`$-torsor under right-multiplication.
#### 4.2.2. The automorphisms $`𝝈_0`$ and $`𝝈_1`$
We define a natural $`\sigma `$-linear automorphisms
$$𝝈_0,𝝈_1:L[X,1/detX]L[X,1/detX],$$
by setting $`𝝈_0X:=X`$ and $`𝝈_1X:=\mathrm{\Phi }X`$. We note that $`𝝈_0`$ restricts to an automorphism of $`R[X,1/detX]`$ for any $`F`$-subalgebra $`R`$ of $`L`$, and that $`𝝈_1`$ induces automorphisms of $`K[X,1/detX]`$ and $`\mathrm{\Sigma }[X,1/detX]`$. We see that
$$𝝈_0𝔮=𝔮,𝝈_1𝔭=𝔭.$$
The first equality is clear since $`𝔮F[X,1/detX]`$. For the second, we note that for $`h(X)K[X,1/detX]`$, we have $`𝝈_1(h)(X)=𝝈_0(h)(\mathrm{\Phi }X)`$, and so
$$𝝈_1(h)(\mathrm{\Psi })=𝝈_0(h)(\mathrm{\Phi }\mathrm{\Psi })=𝝈_0(h)(\sigma (\mathrm{\Psi }))=\sigma (h(\mathrm{\Psi })).$$
Thus, $`𝝈_1𝔭=𝔭`$. This equality implies further that $`\nu 𝝈_1=\sigma \nu `$. The following lemma provides a correspondence between the contraction and extension of ideals in $`L[X,1/detX]`$. See also \[26, Lem. 1.23\].
###### Lemma 4.2.3.
The functions between sets of ideals,
$$\begin{array}{ccc}\{𝔞F[X,1/detX]\}& & \{𝔟L[X,1/detX]𝝈_0𝔟=𝔟\},\\ 𝔞& & (𝔞)\\ 𝔟F[X,1/detX]& & 𝔟\end{array}$$
are bijections.
###### Proof.
Since $`𝝈_0`$ is trivial on $`F[X,1/detX]L[X,1/detX]`$, these maps are well-defined. One knows already that $`(𝔞)F[X,1/detX]=𝔞`$ for all ideals $`𝔞F[X,1/detX]`$ (see \[34, §VII.11\]). Now let $`𝔟L[X,1/detX]`$ be an ideal with $`𝝈_0𝔟=𝔟`$, and let $`𝔞:=𝔟F[X,1/detX]`$. Letting $`\{g_i\}_{iI}`$ be an $`F`$-basis of $`F[X,1/detX]`$, we have that $`\{g_i\}_{iI}`$ is an $`L`$-basis of $`L[X,1/detX]`$. For $`h𝔟`$ we write $`h=b_ig_i`$, $`b_iL`$, and we let $`l(h)`$ be the number of $`iI`$ for which $`b_i0`$. We show that $`h(𝔞)`$ by induction on $`l(h)`$. If $`l(h)=0`$ the result is clear. If $`l(h)=1`$, then $`h=bg`$ for some $`bL^\times `$ and $`g\{g_i\}`$. Moreover, then $`g𝔞`$. Now suppose that $`l(h)>1`$. By multiplying by an element of $`L`$ we can assume that $`b_{i_1}=1`$ and that $`b_{i_2}LF`$ for some $`i_1,i_2I`$. (If all $`b_iF`$, then $`h𝔞`$.) One sees that
$$l(𝝈_0hh)<l(h),$$
and since $`𝝈_0𝔟=𝔟`$, we have $`𝝈_0hh𝔟`$. Therefore, $`𝝈_0(h)h(𝔞)`$. Similarly, $`𝝈_0(b_{i_2}^1h)b_{i_2}^1h(𝔞)`$. However,
$$(\sigma (b_{i_2}^1)b_{i_2}^1)h=\left(𝝈_0(b_{i_2}^1h)b_{i_2}^1h\right)\sigma (b_{i_2}^1)(𝝈_0hh).$$
The left-hand side is non-zero, and the right-hand side is in $`(𝔞)`$. Therefore $`h(𝔞)`$. ∎
###### Proposition 4.2.4.
Define a morphism of affine $`L`$-schemes $`\varphi :=Z_L\mathrm{GL}_{r/L}`$ so that on points $`u\mathrm{\Psi }^1u`$ for $`uZ(\overline{L})`$. Then $`\varphi `$ factors through an isomorphism $`\varphi ^{}:Z_L\mathrm{\Gamma }_L`$ of affine $`L`$-schemes.
###### Proof.
For commutative rings $`RS`$ and for any ideal $`I`$ in $`R[X,1/detX]`$, we let $`I_S`$ denote its extension to $`S[X,1/detX]`$. Now the ideal $`𝔭K[X,1/detX]`$ is the defining ideal of the $`K`$-scheme $`Z`$, and $`𝔮F[X,1/detX]`$ is the defining ideal of of the $`F`$-scheme $`\mathrm{\Gamma }`$. If we set
$$\alpha :L[X,1/detX]L[X,1/detX],$$
to be the $`L`$-algebra homomorphism determined by setting $`\alpha (X)=\mathrm{\Psi }^1X`$, then the map
$$\overline{\alpha }:L[X,1/detX]L[X,1/detX]/𝔭_L,$$
induced by $`\alpha `$, is the map $`\varphi `$ on the level of coordinate rings. It then suffices to prove that $`𝔮_L=\alpha ^1𝔭_L`$.
As noted in §4.2.2, we have that $`𝝈_0𝔮_L=𝔮_L`$ and $`𝝈_1𝔭_L=𝔭_L`$. Furthermore,
$$𝝈_1\alpha X=𝝈_1(\mathrm{\Psi }^1X)=(\sigma \mathrm{\Psi })^1(𝝈_1X)=\mathrm{\Psi }^1X=\alpha 𝝈_0X,$$
and so $`𝝈_1\alpha =\alpha 𝝈_0`$, which implies that
$$𝝈_0\alpha ^1𝔭_L=\alpha ^1𝔭_L.$$
By Lemma 4.2.3, it follows that $`\alpha ^1𝔭_L`$ is generated by $`\alpha ^1𝔭_LF[X,1/detX]`$.
Now we regard $`L_KL`$ as an $`L`$-algebra through the map $`ff1`$. If we let $`\stackrel{~}{\mu }:L[X,1/detX]L_KL`$ be the unique $`L`$-algebra homomorphism such that $`\stackrel{~}{\mu }X=\mathrm{\Psi }_2`$, then we note that the composition
$$F[X,1/detX]\stackrel{\alpha }{}L[X,1/detX]\stackrel{\stackrel{~}{\mu }}{}L_KL$$
is in fact $`\mu `$. Since $`L`$ is a field, the map $`L[X,1/detX]/𝔭_LL_KL`$ induced by $`\stackrel{~}{\mu }`$ is injective. Therefore,
$$𝔮=\alpha ^1𝔭_LF[X,1/detX],$$
and by our argument in the previous paragraph, $`𝔮_L=\alpha ^1𝔭_L`$. ∎
###### Corollary 4.2.5.
The ideal $`𝔭K[X,1/detX]`$ is maximal among proper $`𝛔_1`$-invariant ideals.
###### Proof.
Let $`𝔪𝔭`$ be a proper ideal of $`K[X,1/detX]`$ such that $`𝝈_1𝔪𝔪`$. Because $`K[X,1/detX]`$ is noetherian, it follows that $`𝝈_1𝔪=𝔪`$. Now $`\alpha ^1𝔪_L\alpha ^1𝔭_L=𝔮_L`$, and we see easily that $`𝝈_0\alpha ^1𝔪_L=\alpha ^1𝔪_L`$. Therefore, by Lemma 4.2.3,
$$\alpha ^1𝔪_L=(\alpha ^1𝔪_LF[X,1/detX])_L.$$
Let $`𝔞F[X,1/detX]`$ be a maximal ideal that contains $`\alpha ^1𝔪_LF[X,1/detX]`$, and let $`E:=F[X,1/detX]/𝔞`$, which is a finite extension of $`F`$. By Lemma 4.2.3, we see that $`𝔞=𝔞_LF[X,1/detX]`$, and it follows that there is an isomorphism $`\beta :L[X,1/detX]/𝔞_L\stackrel{}{}L_FE`$. Now if we consider the maps
$$\mathrm{\Pi }:L[X,1/detX]\stackrel{\alpha ^1}{}L[X,1/detX]\stackrel{\beta }{}L_FE,$$
we see that $`𝔪_L\mathrm{ker}\mathrm{\Pi }`$. If we let $`\pi :K[X,1/detX]L_FE`$ be the restriction of $`\mathrm{\Pi }`$, then easily $`𝔪\mathrm{ker}\pi `$ and $`\mathrm{ker}\pi `$ is a proper ideal. Moreover, since $`\alpha ^1𝝈_1=𝝈_0\alpha ^1`$, it follows that $`\mathrm{ker}\pi `$ is a $`𝝈_1`$-invariant ideal of $`K[X,1/detX]`$. Therefore, the maximality of $`𝔪`$ implies that $`𝔪=\mathrm{ker}\pi `$.
Now let $`\mathrm{\Psi }^{}\mathrm{GL}_r(L_FE)`$ be defined by $`\mathrm{\Psi }_{ij}^{}=\pi (X_{ij})`$. The automorphism $`\sigma `$ on $`L`$ extends to an automorphism of $`L_FE`$ by acting by the identity on $`E`$, and it is easily seen that $`(L_FE)^\sigma =E`$. In this way $`\sigma (\mathrm{\Psi }^{})=\mathrm{\Phi }\mathrm{\Psi }^{}`$, and this implies that the matrix $`\delta :=(\mathrm{\Psi }^{})^1\mathrm{\Psi }\mathrm{GL}_r(E)`$. Now $`\delta `$ induces an automorphism on $`(K_FE)[X,1/detX]`$ via
$$\delta h(X):=h(X\delta ).$$
If we extend $`\pi `$ to $`\pi ^{}:(K_FE)[X,1/detX]L_FE`$ by the identity on $`E`$, then we see that we have the extended ideals
$$𝔭_{K_FE}=𝔭_FE𝔪_FE\mathrm{ker}\pi ^{}=\delta (𝔭_FE).$$
But $`(K_FE)[X,1/detX]`$ is a noetherian ring, and so $`𝔭_FE\delta (𝔭_FE)`$ implies that $`𝔭_FE=\delta (𝔭_FE)`$. Thus,
$$𝔭_FE=𝔪_FE=\delta (𝔭_FE).$$
Now $`(K_FE)[X,1/detX]`$ is a free $`K[X,1/detX]`$-module, since $`E`$ is a vector space over $`F`$, and is therefore faithfully flat over $`K[X,1/detX]`$. So by intersecting with $`K[X,1/detX]`$ we see that $`𝔭=𝔪`$ (see \[22, Thm. 7.5\]). ∎
#### 4.2.6. Contracted ideals of $`\mathrm{\Sigma }[X,1/detX]`$
The following lemma is a companion to Lemma 4.2.3, and relies on the preceding corollary. See also \[25, Lem. 1.11\].
###### Lemma 4.2.7.
Let $`𝔟\mathrm{\Sigma }[X,1/detX]`$ be an ideal that is $`𝛔_0`$-invariant. Then $`𝔟`$ is generated by $`𝔟F[X,1/detX]`$.
###### Proof.
Let $`𝔞:=𝔟F[X,1/detX]`$. Let $`\{g_i\}_{iI}`$ be an $`F`$-basis for $`F[X,1/detX]`$ such that $`I=I_𝔞I_1`$, where $`\{g_i\}_{iI_𝔞}`$ is an $`F`$-basis for $`𝔞`$. Choose a subset $`JI_1`$ minimal so that $`𝔟_{iJ}\mathrm{\Sigma }g_i`$ contains a non-zero element of $`𝔟/(𝔞)`$. Pick $`jJ`$, and let
$$𝔪:=\{b\mathrm{\Sigma }_{iJ}b_ig_i𝔟/(𝔞),b_j=b\}$$
Since $`𝝈_0𝔟=𝔟`$, and since each $`g_i`$ is fixed by $`𝝈_0`$, it follows that $`𝔪`$ is a non-zero $`\sigma `$-invariant ideal of $`\mathrm{\Sigma }`$. However, by Corollary 4.2.5, $`\mathrm{\Sigma }`$ has no $`\sigma `$-invariant ideals other than $`\{0\}`$ and $`\mathrm{\Sigma }`$. Thus, $`𝔪=\mathrm{\Sigma }`$. Therefore there exists $`h𝔟`$ so that $`h=_{iJ}b_ig_imod(𝔞)`$ and $`b_j=1`$. Now $`𝝈_0(h)h`$ is supported on a proper subset of $`J`$ modulo $`(𝔞)`$, and so it must be $`0`$ modulo $`(𝔞)`$ by the minimality of $`J`$. Therefore each $`b_iF`$, and thus $`_{iJ}b_ig_i(𝔞)`$, which is a contradiction. ∎
###### Proposition 4.2.8.
Define a morphism of affine $`K`$-schemes $`\psi :Z\times ZZ\times \mathrm{GL}_{r/K}`$ so that on points $`(u,v)(u,u^1v)`$ for $`u`$, $`vZ(\overline{K})`$. Then $`\psi `$ factors through an isomorphism $`Z\times ZZ\times \mathrm{\Gamma }_K`$ of affine $`K`$-schemes.
###### Proof.
Again we work on the level of coordinate rings and maintain conventions and definitions in the proof of Proposition 4.2.4. The ring $`\mathrm{\Sigma }L`$ is isomorphic to the coordinate ring of $`Z`$ over $`K`$. Likewise, the ring $`\mathrm{\Sigma }[X,1/detX]/𝔭_\mathrm{\Sigma }`$ is the coordinate ring of $`Z\times Z`$, and $`\mathrm{\Sigma }[X,1/detX]/𝔮_\mathrm{\Sigma }`$ is the coordinate ring of $`Z\times \mathrm{\Gamma }_K`$. The $`L`$-algebra automorphism $`\alpha `$ in the proof of the previous proposition restricts to an automorphism of $`\mathrm{\Sigma }[X,1/detX]`$, and in this way the homomorphism
$$\overline{\alpha }:\mathrm{\Sigma }[X,1/detX]\mathrm{\Sigma }[X,1/detX]/𝔭_\mathrm{\Sigma }$$
induced by $`\alpha `$ represents the morphism $`\psi `$ of affine $`K`$-schemes. We then need to show that $`𝔮_\mathrm{\Sigma }=\alpha ^1𝔭_\mathrm{\Sigma }`$.
Let $`𝔞=\alpha ^1𝔭_\mathrm{\Sigma }F[X,1/detX]`$. By Lemma 4.2.7,
$$\alpha ^1𝔭_\mathrm{\Sigma }=𝔞_\mathrm{\Sigma }.$$
Then as in the proof of Proposition 4.2.4, $`𝔮_L=\alpha ^1𝔭_L`$ and both are now equal to $`𝔞_L`$. Since $`𝔮`$ and $`𝔞`$ are both ideals in $`F[X,1/detX]`$ and $`F`$ and $`L`$ are fields, it follows that $`𝔮=𝔞`$. ∎
###### Lemma 4.2.9.
Let $`G`$ be a group, and let $`A`$ and $`B`$ be subsets of $`G`$, $`A`$ non-empty, such that the map
$$(u,v)(u,u^1v):A\times AA\times G$$
factors through a bijection $`\varphi :A\times AA\times B`$. Then $`B`$ is a subgroup of $`G`$ and $`A`$ is stable under right-multiplication by elements of $`B`$. Moreover, under the action of $`B`$ by right-multiplication, $`A`$ becomes a principal homogeneous space for $`B`$.
###### Proof.
This is a simple exercise. ∎
#### 4.2.10. The Galois group $`\mathrm{\Gamma }`$
The previous propositions and lemmas culminate in the following theorem, saying that $`\mathrm{\Gamma }`$ is in fact an affine group scheme. We call the group $`\mathrm{\Gamma }`$, or $`\mathrm{\Gamma }_\mathrm{\Psi }`$ if we wish to recall the dependence on $`\mathrm{\Psi }`$, the *Galois group of the system $`\sigma (\mathrm{\Psi })=\mathrm{\Phi }\mathrm{\Psi }`$*.
###### Theorem 4.2.11.
Let $`(F,K,L)`$ be a $`\sigma `$-admissible triple for an automorphism $`\sigma :LL`$. Suppose we have $`\mathrm{\Phi }GL_r(K)`$ and $`\mathrm{\Psi }GL_r(L)`$ so that $`\sigma (\mathrm{\Psi })=\mathrm{\Phi }\mathrm{\Psi }`$. Then $`\mathrm{\Gamma }:=\mathrm{\Gamma }_\mathrm{\Psi }`$ is a closed $`F`$-subgroup scheme of $`\mathrm{GL}_{r/F}`$, and the closed $`K`$-subscheme $`Z:=Z_\mathrm{\Psi }`$ of $`\mathrm{GL}_{r/K}`$ is stable under right-multiplication by $`\mathrm{\Gamma }_K`$ and is a $`\mathrm{\Gamma }_K`$-torsor.
###### Proof.
Since $`Z(\mathrm{\Sigma })`$ is non-empty, Propositions 4.2.4 and 4.2.8 imply that $`(u,v)(u,u^1v):Z(\mathrm{\Sigma })\times Z(\mathrm{\Sigma })Z(\mathrm{\Sigma })\times \mathrm{\Gamma }(\mathrm{\Sigma })`$ is a bijection. Lemma 4.2.9 and the Yoneda lemma \[32, §1.2–1.4\] imply that $`\mathrm{\Gamma }_\mathrm{\Sigma }`$ is a subgroup of $`\mathrm{GL}_{r/\mathrm{\Sigma }}`$ and that $`Z_\mathrm{\Sigma }`$ is a $`\mathrm{\Gamma }_\mathrm{\Sigma }`$-torsor. Since the inclusion $`F\mathrm{\Sigma }`$ is faithfully flat, we see that $`\mathrm{\Gamma }`$ is a closed $`F`$-subgroup scheme of $`\mathrm{GL}_{r/F}`$ by flat descent \[32, §17.1–17.3\]. Similarly, since the inclusion $`K\mathrm{\Sigma }`$ is faithfully flat, $`Z`$ admits the structure of a $`\mathrm{\Gamma }_K`$-torsor. ∎
### 4.3. Criterion for smoothness
We continue with the notation of the previous section, and in particular have fixed a $`\sigma `$-admissible triple $`(F,K,L)`$ together with $`\mathrm{\Phi }\mathrm{GL}_r(K)`$, $`\mathrm{\Psi }\mathrm{GL}_r(L)`$ satisfying $`\sigma (\mathrm{\Psi })=\mathrm{\Phi }\mathrm{\Psi }`$. In this section, we explore when $`\mathrm{\Gamma }`$ is smooth over $`\overline{F}`$, that is, when the coordinate ring of $`\mathrm{\Gamma }_{\overline{F}}`$ is reduced.
###### Theorem 4.3.1.
Suppose $`K`$ is algebraically closed in the fraction field $`\mathrm{\Lambda }`$ of $`\mathrm{\Sigma }`$. Then
1. The $`K`$-scheme $`Z`$ is absolutely irreducible and is smooth over $`\overline{K}`$.
2. The $`F`$-scheme $`\mathrm{\Gamma }`$ is absolutely irreducible and is smooth over $`\overline{F}`$.
3. The dimension of $`\mathrm{\Gamma }`$ over $`F`$ is equal to the transcendence degree of $`\mathrm{\Lambda }`$ over $`K`$.
###### Proof.
The ideal $`𝔭K[X,1/detX]`$ is prime. The field $`\mathrm{\Lambda }`$ is separable over $`K`$, since it is a subfield of $`L`$. That $`K`$ is algebraically closed in $`\mathrm{\Lambda }`$ then implies that $`𝔭_{\overline{K}}\overline{K}[X,1/detX]`$ is prime \[34, VII.11, Thm. 39\]. Thus $`Z`$ is smooth over $`\overline{K}`$. Because $`\mathrm{\Gamma }_{\overline{K}}Z_{\overline{K}}`$, $`\mathrm{\Gamma }`$ must be smooth over $`\overline{F}`$. By construction the transcendence degree of $`\mathrm{\Lambda }`$ over $`K`$ is equal to the dimension of $`Z`$, which is equal to the dimension of $`\mathrm{\Gamma }`$. ∎
#### 4.3.2. The case $`(𝔽_q(t),\overline{k}(t),𝕃)`$
This case is of particular interest to our applications to $`t`$-motives. It turns out that in this case, all Galois groups are smooth, via the following proposition. We continue with our usual notation.
###### Proposition 4.3.3.
Suppose $`(F,K,L)=(𝔽_q(t),\overline{k}(t),𝕃)`$, and suppose that $`\mathrm{\Phi }\mathrm{GL}_r(\overline{k}(t))`$ and $`\mathrm{\Psi }\mathrm{GL}_r(𝕃)`$ satisfy $`\mathrm{\Psi }^{(1)}=\mathrm{\Phi }\mathrm{\Psi }`$. Then $`\overline{k}(t)`$ is algebraically closed in $`\mathrm{\Lambda }=\overline{k}(t)(\mathrm{\Psi })`$.
###### Proof.
Let $`f\mathrm{\Lambda }\overline{k(t)}`$, and consider the field $`H:=\overline{k}(t;f^{(i)}:i)`$ obtained by adjoining all of the twists of $`f`$ to $`\overline{k}(t)`$. Each $`f^{(i)}`$ is algebraic over $`\overline{k}(t)`$, and so $`H/\overline{k}(t)`$ is algebraic. Since $`\mathrm{\Lambda }`$ is finitely generated as a field over $`\overline{k}(t)`$, so is $`H`$. Thus $`[H:\overline{k}(t)]<\mathrm{}`$. Furthermore, $`H`$ is invariant under $`\sigma `$ and $`\sigma ^1`$.
The field $`H`$ is the function field of a smooth projective curve $`X`$ over $`\overline{k}`$, and the inclusion $`\overline{k}(t)H`$ provides a surjective morphism $`X_{\overline{k}}^1`$ over $`\overline{k}`$. Now $`\sigma :HH`$ induces an automorphism $`\tau :XX`$ as a scheme over $`𝔽_q`$. Because $`\sigma `$ leaves the integral closure of $`\overline{k}[t]`$ in $`H`$ invariant, the points $`\mathrm{}_1,\mathrm{},\mathrm{}_d`$ in $`X`$ above the point $`\mathrm{}`$ in $`_{\overline{k}}^1`$ are permuted by $`\sigma `$. Thus we can construct an effective divisor $`I`$ of $`X`$ such that $`\tau (I)=I`$ and $`\mathrm{Supp}(I)=\{\mathrm{}_1,\mathrm{},\mathrm{}_d\}`$. Now for $`N1`$ sufficiently large, the field $`H`$ is generated over $`\overline{k}(t)`$ by the functions in the finite dimensional $`\overline{k}`$-vector space
$$S:=\mathrm{\Gamma }(X,NI)H.$$
By our assumptions on $`I`$, this space is invariant under $`\sigma `$ and $`\sigma ^1`$. If the entries of $`𝐟:=[f_1,\mathrm{},f_m]^{\text{tr}}`$ form a $`\overline{k}`$-basis for $`S`$, then there is a matrix $`A\mathrm{GL}_m(\overline{k})`$ so that $`\sigma (𝐟)=A𝐟`$. If $`𝐠\mathrm{Mat}_{m\times 1}(S)`$ and $`𝐠=B𝐟`$ for some $`B\mathrm{GL}_m(\overline{k})`$, then
$$\sigma (𝐠)=B^{(1)}AB^1𝐠.$$
By the theory of Lang isogenies , we can pick a $`B\mathrm{GL}_m(\overline{k})`$ so that
$$B^1B^{(1)}=A^{(1)},$$
and if we let $`𝐠:=B𝐟`$, then $`\sigma (𝐠)=𝐠`$. Thus $`S`$ contains a $`\overline{k}`$-basis $`𝐠`$ that is fixed by $`\sigma `$, and $`H=\overline{k}(t,𝐠)`$. Let $`g`$ be an entry of $`𝐠`$. Then $`g\overline{k(t)}𝕃^\sigma =𝔽_q(t)`$. Thus $`[H:\overline{k}(t)]=1`$. ∎
###### Corollary 4.3.4.
Let $`(F,K,L)=(𝔽_q(t),\overline{k}(t),𝕃)`$, and suppose that $`\mathrm{\Phi }\mathrm{GL}_r(\overline{k}(t))`$ and $`\mathrm{\Psi }\mathrm{GL}_r(𝕃)`$ satisfy $`\sigma (\mathrm{\Psi })=\mathrm{\Phi }\mathrm{\Psi }`$. Then the Galois group $`\mathrm{\Gamma }`$ of $`\mathrm{\Psi }`$ is smooth over $`\overline{𝔽_q(t)}`$.
### 4.4. The Galois action
In this section we will assume that $`K`$ is algebraically closed in $`\mathrm{\Lambda }`$, and so in particular by Theorem 4.3.1, $`\mathrm{\Gamma }_{\overline{F}}`$ and $`Z_{\overline{K}}`$ are reduced and irreducible.
#### 4.4.1. $`\sigma `$-automorphisms of $`\mathrm{\Sigma }`$ and $`\mathrm{\Lambda }`$
Let $`\mathrm{Aut}_\sigma (\mathrm{\Sigma }/K)`$ denote the group of automorphisms of $`\mathrm{\Sigma }`$ over $`K`$ that commute with $`\sigma `$. Similarly we define $`\mathrm{Aut}_\sigma (\mathrm{\Lambda }/K)`$. In fact, it is true that
$$\mathrm{Aut}_\sigma (\mathrm{\Sigma }/K)=\mathrm{Aut}_\sigma (\mathrm{\Lambda }/K).$$
Indeed every automorphism $`\xi \mathrm{Aut}_\sigma (\mathrm{\Sigma }/K)`$ extends uniquely to an automorphism in $`\mathrm{Aut}_\sigma (\mathrm{\Lambda }/K)`$. On the other hand, if $`\eta \mathrm{Aut}_\sigma (\mathrm{\Lambda }/K)`$, then as matrices in $`\mathrm{Mat}_r(L)`$, $`\sigma (\eta (\mathrm{\Psi }))=\eta (\sigma (\mathrm{\Psi }))=\eta (\mathrm{\Phi }\mathrm{\Psi })=\mathrm{\Phi }\eta (\mathrm{\Psi })`$. Thus, $`\eta (\mathrm{\Psi })=\mathrm{\Psi }\gamma `$ for some $`\gamma \mathrm{GL}_r(F)`$, and so $`\eta (\mathrm{\Psi })\mathrm{Mat}_r(\mathrm{\Sigma })`$. Therefore $`\eta `$ restricted to $`\mathrm{\Sigma }`$ takes values in $`\mathrm{\Sigma }`$.
#### 4.4.2. The action of $`\mathrm{\Gamma }(F)`$
For $`\gamma \mathrm{\Gamma }(F)`$, we have an automorphism of $`K`$-schemes $`\gamma :ZZ`$ defined by right multiplication by $`\gamma `$. On the level of coordinate rings, the induced map is
$$\gamma =(h(X)h(X\gamma )):\mathrm{\Sigma }\mathrm{\Sigma },$$
which is a $`K`$-linear automorphism that commutes with the action of $`\sigma `$. Thus we have a group homomorphism,
$$\kappa :\mathrm{\Gamma }(F)\mathrm{Aut}_\sigma (\mathrm{\Lambda }/K),$$
which is easily seen to be injective. Now if $`\delta \mathrm{Aut}_\sigma (\mathrm{\Lambda }/K)`$, then $`\delta `$ induces an automorphism of the non-empty $`Z(\mathrm{\Lambda })`$ that is right-multiplication by an element of $`\gamma \mathrm{\Gamma }(\mathrm{\Lambda })`$. That $`\delta `$ commutes with $`\sigma `$ implies that $`\gamma \mathrm{\Gamma }(\mathrm{\Lambda }^\sigma )=\mathrm{\Gamma }(F)`$. Thus $`\kappa `$ is an isomorphism.
#### 4.4.3. Base extensions
Given our $`\sigma `$-admissible triple $`(F,K,L)`$, we choose an extension of $`\sigma `$ to an automorphism of $`\overline{L}`$. Then $`\overline{L}^\sigma `$ is an algebraic extension of $`F`$. Indeed, the monic irreducible polynomial of any $`h\overline{L}^\sigma `$ over $`L`$ must have coefficients in $`L^\sigma =F`$. Thus if we let $`𝐅=\overline{L}^\sigma `$, then $`(𝐅,\overline{K},\overline{L})`$ is a $`\sigma `$-admissible triple. The Galois group $`\mathrm{\Gamma }^{}`$ defined by the system $`\sigma (\mathrm{\Psi })=\mathrm{\Phi }\mathrm{\Psi }`$ defined with respect to $`(𝐅,\overline{K},\overline{L})`$ is seen to be $`\mathrm{\Gamma }_𝐅`$ by Propositions 4.2.4 and 4.2.8. If we let $`\stackrel{~}{\mathrm{\Sigma }}`$ be the coordinate of $`Z_{\overline{K}}`$ and $`\stackrel{~}{\mathrm{\Lambda }}`$ be its fraction field, then we see that
$$\mathrm{\Gamma }(𝐅)\mathrm{Aut}_\sigma (\stackrel{~}{\mathrm{\Sigma }}/\overline{K})=\mathrm{Aut}_\sigma (\stackrel{~}{\mathrm{\Lambda }}/\overline{K}).$$
Furthermore, for $`n1`$, let $`𝐅_n=\overline{L}^{\sigma ^n}`$, and suppose that $`(𝐅_n,\overline{K},\overline{L})`$ is $`\sigma ^n`$-admissible. Then $`\mathrm{\Psi }`$ is a fundamental matrix for $`\mathrm{\Phi }_n:=\sigma ^{n1}(\mathrm{\Phi })\mathrm{}\sigma (\mathrm{\Phi })\mathrm{\Phi }`$. Again by Propositions 4.2.4 and 4.2.8, we see that the Galois group of this system of equations is $`\mathrm{\Gamma }_{𝐅_n}`$. And thus,
$$\mathrm{\Gamma }(𝐅_n)\mathrm{Aut}_{\sigma ^n}(\stackrel{~}{\mathrm{\Sigma }}/\overline{K})=\mathrm{Aut}_{\sigma ^n}(\stackrel{~}{\mathrm{\Lambda }}/\overline{K}).$$
#### 4.4.4. Galois action for $`\mathrm{\Gamma }(\overline{F})`$
Continuing with the notation of the previous paragraphs, suppose that $`\overline{F}=𝐅_n`$. Then every element of $`\mathrm{\Gamma }(\overline{F})`$ induces an automorphism of $`\stackrel{~}{\mathrm{\Lambda }}/\overline{K}`$ that commutes with $`\sigma ^n`$ for all $`n0`$. In this case, we will call this the induced action of $`\mathrm{\Gamma }(\overline{F})`$ on $`\stackrel{~}{\mathrm{\Lambda }}`$.
#### 4.4.5. The case $`(𝔽_q(t),\overline{k}(t),𝕃)`$
It is worth pointing out that the situation is quite nice in our usual setting. For $`n1`$, the triple $`(𝔽_{q^n}(t),\overline{k}(t),𝕃)`$ is $`\sigma ^n`$-admissible. As in §2.2.5, there is a canonical extension of $`\sigma `$ to $`𝕂t\overline{𝕃}`$. Furthermore, we see that
$$\overline{𝕃}^{\sigma ^n}=\overline{k(t)}^{\sigma ^n}=\overline{𝔽_q(t)}^{\sigma ^n}=:𝐅_n,$$
and so $`(𝐅_n,\overline{k(t)},\overline{𝕃})`$ is a $`\sigma ^n`$-admissible triple. Every element of $`\overline{𝔽_q(t)}`$ is fixed by some power of $`\sigma `$, and so
$$\overline{𝔽_q(t)}=\underset{n1}{}𝐅_n.$$
We now return to the general situation, but it is important to note that the following theorem applies to Galois groups in the usual $`(𝔽_q(t),\overline{k}(t),𝕃)`$ setting.
###### Theorem 4.4.6.
Let $`\mathrm{\Phi }\mathrm{GL}_r(K)`$, and suppose that $`\mathrm{\Psi }\mathrm{GL}_r(L)`$ is a fundamental matrix for $`\mathrm{\Phi }`$. Assume that $`K`$ is algebraically closed in $`\mathrm{\Lambda }=K(\mathrm{\Psi })`$. Fix an extension of $`\sigma `$ to $`\overline{L}`$, and let $`𝐅_n:=\overline{L}^{\sigma ^n}`$. Suppose that $`(𝐅_n,\overline{K},\overline{L})`$ is $`\sigma ^n`$-admissible for each $`n1`$, and suppose that $`\overline{F}=𝐅_n`$. Let $`\stackrel{~}{\mathrm{\Sigma }}`$ be the coordinate ring of $`Z_{\overline{K}}`$ and let $`\stackrel{~}{\mathrm{\Lambda }}`$ be its fraction field, both considered subrings of $`\overline{L}`$.
1. The subfield of $`\stackrel{~}{\mathrm{\Lambda }}`$ fixed by $`\mathrm{\Gamma }(\overline{F})`$ is $`\overline{K}`$.
2. The elements of $`\mathrm{\Lambda }`$ fixed by $`\mathrm{\Gamma }(\overline{F})`$ are precisely $`K`$.
###### Proof.
See \[25, Lem. 1.28\]. Suppose $`f\stackrel{~}{\mathrm{\Lambda }}`$ is fixed by $`\mathrm{\Gamma }(\overline{F})`$. We consider $`f\stackrel{~}{\mathrm{\Lambda }}`$ to be a function $`f:Z_{\overline{K}}_{\overline{K}}^1`$. For $`i=1,2`$, we consider the two maps of $`\overline{K}`$-schemes
$$g_i:Z_{\overline{K}}\times \mathrm{\Gamma }_{\overline{K}}Z_{\overline{K}}\times Z_{\overline{K}}\stackrel{\pi _i}{}Z_{\overline{K}}\stackrel{f}{}_{\overline{K}}^1,$$
where $`\pi _i`$ is the $`i`$-th projection. Because $`f`$ is $`\mathrm{\Gamma }(\overline{F})`$-invariant and because $`\mathrm{\Gamma }(\overline{F})`$ is dense in $`\mathrm{\Gamma }_{\overline{K}}`$ since $`\mathrm{\Gamma }`$ is smooth over $`\overline{F}`$, we must have $`g_1=g_2`$. Therefore, $`f\pi _1=f\pi _2`$, which implies that $`f`$ is constant. This proves part (a). Part (b) follows from part (a) and the assumption that $`K`$ is algebraically closed in $`\mathrm{\Lambda }`$. ∎
#### 4.4.7. Remark
If $`\mathrm{\Gamma }(F)`$ is Zariski dense in $`\mathrm{\Gamma }`$, then it follows that $`\mathrm{\Lambda }^{\mathrm{\Gamma }(F)}=\overline{k}(t)`$.
### 4.5. The group $`\mathrm{\Gamma }`$ and $`t`$-motives
Given a $`t`$-motive $`M`$, we defined the Galois group $`\mathrm{\Gamma }_M`$ of $`M`$ in §3.5.2. Associated to $`M`$ we can also choose a matrix $`\mathrm{\Phi }\mathrm{GL}_r(\overline{k}(t))`$ that represents multiplication by $`𝝈`$ on $`M`$. Let $`\mathrm{\Psi }\mathrm{GL}_r(𝕃)`$ be a rigid analytic trivialization of $`\mathrm{\Phi }`$. We will show that $`\mathrm{\Gamma }_M`$ is isomorphic to $`\mathrm{\Gamma }:=\mathrm{\Gamma }_\mathrm{\Psi }`$ over $`𝔽_q(t)`$.
#### 4.5.1. $`t`$-motives and $`\sigma `$-semilinear equations
Let $`M`$ be a $`t`$-motive. We fix the following notation throughout this section. Let $`𝐦`$ be a basis for $`M`$, and let $`\mathrm{\Phi }\mathrm{GL}_r(\overline{k}(t))`$ represent multiplication by $`𝝈`$ on $`M`$. We pick a rigid analytic trivialization $`\mathrm{\Psi }\mathrm{GL}_r(𝕃)`$ for $`M`$, which is at the same time a fundamental matrix for $`\mathrm{\Phi }`$.
Let $`M_v^u:=M^u(M^{})^v`$. Because $`𝒯_M`$ is Tannakian, if $`N`$ is any $`t`$-motive in $`𝒯_M`$, then $`N`$ is the subquotient of a direct sum of various $`M_v^u`$, and vice versa. It follows from Propositions 3.3.9(b) and 3.3.11(b) that the entries of a fundamental matrix $`\mathrm{\Psi }_N`$ for $`N`$ are in $`\mathrm{\Sigma }`$, and in fact we can take $`\mathrm{\Psi }_N\mathrm{GL}_s(\mathrm{\Sigma })`$ for some $`s`$.
For an $`𝔽_q(t)`$-algebra $`R`$, we let $`\mathrm{\Sigma }^{(R)}:=R_{𝔽_q(t)}\mathrm{\Sigma }`$.
###### Lemma 4.5.2.
For any $`t`$-motive $`N`$ in $`𝒯_M`$ and $`𝔽_q(t)`$-algebra $`R`$, the natural map,
$$\mathrm{\Sigma }^{(R)}_{𝔽_q(t)}N^\mathrm{B}\mathrm{\Sigma }^{(R)}_{\overline{k}(t)}N$$
is bijective.
###### Proof.
Let $`\kappa `$ be the map defined in the statement of the lemma. Thus as above we can pick a basis $`𝐧`$ for $`N`$ and a rigid analytic trivialization $`\mathrm{\Psi }_N\mathrm{GL}_s(\mathrm{\Sigma })`$ with respect to $`𝐧`$. By Proposition 3.3.9(b), $`\mathrm{\Psi }_N^1𝐧`$ is an $`𝔽_q(t)`$-basis for $`N^\mathrm{B}`$. Now $`1(\mathrm{\Psi }_N^1𝐧)`$ is a $`\mathrm{\Sigma }^{(R)}`$-basis of $`\mathrm{\Sigma }^{(R)}_{𝔽_q(t)}N^\mathrm{B}`$ (here and elsewhere $`1A`$ for a matrix $`A`$ is the matrix of the same dimension whose entries are each tensored by 1 on the left). If $`𝐟\mathrm{Mat}_{1\times s}(\mathrm{\Sigma }^{(R)})`$, then
$$\kappa \left((𝐟1)(1(\mathrm{\Psi }_N^1𝐧))\right)=(𝐟\mathrm{\Psi }_N^11)(1𝐧).$$
The entries of $`(\mathrm{\Psi }_N^11)(1𝐧)`$ are in the image of $`\kappa `$, and
$$\mathrm{\Psi }_N(\mathrm{\Psi }_N^11)(1𝐧)=1𝐧.$$
Thus $`\kappa `$ is surjective. Since $`\mathrm{\Psi }_N\mathrm{GL}_s(\mathrm{\Sigma })`$, the map $`\kappa `$ is bijective. ∎
###### Theorem 4.5.3.
Let $`M`$ be a $`t`$-motive, and let $`N`$ be a $`t`$-motive in $`𝒯_M`$. If we consider $`N^\mathrm{B}`$ to be an algebraic group over $`𝔽_q(t)`$, then there is a natural representation
$$\xi _N:\mathrm{\Gamma }\mathrm{GL}(N^\mathrm{B})$$
over $`𝔽_q(t)`$ that is functorial in $`N`$.
###### Proof.
Since every $`t`$-motive $`N`$ in $`𝒯_M`$ is constructed from $`M`$ via tensor products, duals, and subquotients, to define this representation it suffices to define it on $`M^\mathrm{B}`$ itself. Functoriality in $`N`$ will be automatic.
To define the representation on $`M^\mathrm{B}`$, it suffices by the Yoneda lemma \[32, §1.2–1.4\] to define a representation
$$\xi _M^{(R)}:\mathrm{\Gamma }(R)\mathrm{GL}(R_{𝔽_q(t)}M^\mathrm{B})$$
for every $`𝔽_q(t)`$-algebra $`R`$ and show that it is functorial in $`R`$. Let $`R`$ be an $`𝔽_q(t)`$-algebra, and let $`\gamma \mathrm{\Gamma }(R)`$. Define
$$\mathrm{\Xi }^{(R)}(\gamma ):=\gamma 1=(h(\mathrm{\Psi })mh(\mathrm{\Psi }\gamma )m):\mathrm{\Sigma }^{(R)}_{\overline{k}(t)}M\mathrm{\Sigma }^{(R)}_{\overline{k}(t)}M,$$
which is an isomorphism of $`\overline{k}(t)`$-vector spaces. Now by Lemma 4.5.2, $`R_{𝔽_q(t)}M^\mathrm{B}`$ spans $`\mathrm{\Sigma }^{(R)}_{\overline{k}(t)}M`$ as a $`\mathrm{\Sigma }^{(R)}`$-module. Let $`\xi ^{(R)}(\gamma )`$ be the restriction of $`\mathrm{\Xi }^{(R)}(\gamma )`$ to $`R_{𝔽_q(t)}M^\mathrm{B}`$.
We claim that the image of $`\xi ^{(R)}(\gamma )`$ is $`R_{𝔽_q(t)}M^\mathrm{B}`$. Indeed, since $`\mathrm{\Psi }^1𝐦`$ forms an $`𝔽_q(t)`$-basis of $`M^\mathrm{B}`$ by Proposition 3.3.9(b), for $`𝐟\mathrm{Mat}_{1\times r}(R)`$, we have
$$\xi ^{(R)}(\gamma ):𝐟(1\mathrm{\Psi }^1)𝐦𝐟\gamma ^1(1\mathrm{\Psi }^1)𝐦.$$
Thus $`\xi ^{(R)}(\gamma )`$ is an $`R`$-linear automorphism of $`R_{𝔽_q(t)}M^\mathrm{B}`$. It is straightforward to check that this is construction is functorial in $`R`$, and so we have defined a homomorphism $`\xi _M:\mathrm{\Gamma }\mathrm{GL}(M^\mathrm{B})`$. ∎
###### Corollary 4.5.4.
Let $`M`$ be a $`t`$-motive. The representation $`\xi _M:\mathrm{\Gamma }\mathrm{GL}(M^\mathrm{B})`$ is faithful.
###### Proof.
As defined in the the proof of the previous theorem we see easily that $`\xi ^{(R)}:\mathrm{\Gamma }(R)\mathrm{GL}(R_{𝔽_q(t)}M^\mathrm{B})`$ is injective for all $`𝔽_q(t)`$-algebras $`R`$. ∎
#### 4.5.5. The functor $`\xi _M`$
For a $`t`$-motive $`M`$, if $`\mathrm{\Phi }\mathrm{GL}_r(\overline{k}(t))`$ represents multiplication by $`𝝈`$ on $`M`$ and if $`\mathrm{\Psi }\mathrm{GL}_r(𝕃)`$ is a rigid analytic trivialization of $`\mathrm{\Phi }`$, then Theorem 4.5.3 defines a functor
$$\xi _M:𝒯_M\mathrm{𝐑𝐞𝐩}(\mathrm{\Gamma },𝔽_q(t)).$$
It is straightforward to check that $`\xi _M`$ is a tensor functor. Let
$$\eta _M:\mathrm{𝐑𝐞𝐩}(\mathrm{\Gamma }_M,𝔽_q(t))\stackrel{}{}𝒯_M$$
be the equivalence of categories defined in §3.5.2. Letting
$$F:\mathrm{𝐑𝐞𝐩}(\mathrm{\Gamma },𝔽_q(t))\mathrm{𝐕𝐞𝐜}(𝔽_q(t))$$
be the forgetful functor, we see immediately that $`\omega _M=F\xi _M`$. Thus by \[11, Cor. II.2.9\], there is a unique homomorphism $`\pi _M:\mathrm{\Gamma }\mathrm{\Gamma }_M`$ over $`𝔽_q(t)`$ so that the natural functor $`\tau _M:\mathrm{𝐑𝐞𝐩}(\mathrm{\Gamma }_M,𝔽_q(t))\mathrm{𝐑𝐞𝐩}(\mathrm{\Gamma },𝔽_q(t))`$ induced by $`\pi _M`$ satisfies
$$\xi _M\eta _M=\tau _M.$$
###### Proposition 4.5.6.
Let $`M`$ be a $`t`$-motive. Suppose that $`\mathrm{\Phi }\mathrm{GL}_r(\overline{k}(t))`$ represents multiplication by $`𝛔`$ on $`M`$ and that $`\mathrm{\Psi }\mathrm{GL}_r(𝕃)`$ is a rigid analytic trivialization for $`\mathrm{\Phi }`$. Then the functor
$$\xi _M:𝒯_M\mathrm{𝐑𝐞𝐩}(\mathrm{\Gamma },𝔽_q(t))$$
is fully faithful.
###### Proof.
For any $`t`$-motives $`N`$ and $`P`$ in $`𝒯_M`$, there is a natural isomorphism of $`𝔽_q(t)`$-vector spaces, $`\mathrm{Hom}_{𝒯_M}(P,N)\mathrm{Hom}_{𝒯_M}(\mathrm{𝟏},\mathrm{Hom}(P,N))`$. Thus it suffices to prove full faithfulness when $`P=\mathrm{𝟏}`$. Now $`\mathrm{Hom}_{𝒯_M}(\mathrm{𝟏},N)NN^\mathrm{B}=\{nN𝝈n=n\}`$, and this provides an injection $`\mathrm{Hom}_{𝒯_M}(\mathrm{𝟏},N)\mathrm{Hom}_\mathrm{\Gamma }(\mathrm{𝟏}^\mathrm{B},N^\mathrm{B})`$.
Conversely suppose that $`\varphi :\mathrm{𝟏}^\mathrm{B}N^\mathrm{B}`$ is a $`\mathrm{\Gamma }`$-morphism. Pick a $`\overline{k}(t)`$-basis $`𝐧`$ for $`N`$. Then $`\varphi (1)=𝐡(\mathrm{\Psi })𝐧`$ for some $`𝐡(\mathrm{\Psi })\mathrm{Mat}_{1\times s}(\mathrm{\Sigma })`$ by Lemma 4.5.2. Let $`E/𝔽_q(t)`$ be a finite extension of fields. We see that for $`\gamma \mathrm{\Gamma }(E)`$, the action of $`\xi ^{(E)}(\gamma ):=\xi _M^{(E)}(N)(\gamma )`$ on $`E_{𝔽_q(t)}N^\mathrm{B}`$ is simply the restriction of the natural map $`\mathrm{\Xi }^{(E)}(\gamma )=1\gamma :\mathrm{\Sigma }^{(E)}_{\overline{k}(t)}N\mathrm{\Sigma }^{(E)}_{\overline{k}(t)}N`$ to $`E_{𝔽_q(t)}N^\mathrm{B}`$. Since $`\varphi `$ is a $`\mathrm{\Gamma }`$-morphism, it follows that $`\xi ^{(E)}(\gamma )(\varphi (1))=\varphi (1)`$ for all $`\gamma \mathrm{\Gamma }(E)`$. Thus,
$$𝐡(\mathrm{\Psi })𝐧=\varphi (1)=\xi ^{(E)}(\gamma )(\varphi (1))=𝐡(\mathrm{\Psi }\gamma )𝐧,\gamma \mathrm{\Gamma }(E).$$
Because $`𝐧`$ is a $`\mathrm{\Sigma }^{(E)}`$-basis of $`\mathrm{\Sigma }^{(E)}_{\overline{k}(t)}N`$, the entries of $`𝐡(\mathrm{\Psi })`$ must each be fixed by every $`\gamma \mathrm{\Gamma }(E)`$. By varying over all $`E/𝔽_q(t)`$ finite, Theorem 4.4.6 implies that $`𝐡(\mathrm{\Psi })\mathrm{Mat}_{1\times s}(\overline{k}(t))`$. Thus $`\varphi (1)NN^\mathrm{B}`$. ∎
###### Lemma 4.5.7.
Let $`\mathrm{\Phi }\mathrm{GL}_r(\overline{k}(t))`$, and suppose that $`\mathrm{\Psi }\mathrm{GL}_r(𝕃)`$ is a fundamental matrix for $`\mathrm{\Phi }`$. Suppose that $`W\mathrm{\Lambda }^s`$ is a vector subspace over $`\mathrm{\Lambda }`$ such that for every finite extension of fields $`E/𝔽_q(t)`$,
$$\mathrm{\Gamma }(E)(E_{𝔽_q(t)}W)E_{𝔽_q(t)}W.$$
Then $`W`$ has a system of defining equations over $`\overline{k}(t)`$.
###### Proof.
Suppose that $`W`$ has dimension $`sm`$, and let $`A(\mathrm{\Psi })\mathrm{Mat}_{m\times s}(\mathrm{\Lambda })`$ be a coefficient matrix for a system of defining equations for $`W`$. By changing the order of the variables if necessary, we can use Gaussian elimination on $`A(\mathrm{\Psi })`$ to obtain
$$G(\mathrm{\Psi })=[I_m,C(\mathrm{\Psi })],$$
where $`C(\mathrm{\Psi })\mathrm{Mat}_{m\times (sm)}(\mathrm{\Lambda })`$. Both $`A(\mathrm{\Psi })`$ and $`G(\mathrm{\Psi })`$ provide coefficient matrices for equations for $`W`$, and so it suffices to show that $`C(\mathrm{\Psi })`$ has entries in $`\overline{k}(t)`$.
Let $`E/𝔽_q(t)`$ be a finite extension of fields. Since $`E_{𝔽_q(t)}W`$ is invariant under $`\mathrm{\Gamma }(E)`$, it follows that, for every $`\gamma \mathrm{\Gamma }(E)`$, the matrix $`G(\mathrm{\Psi }\gamma ^1)`$ is also the coefficient matrix of a defining set of equations for $`E_{𝔽_q(t)}W`$. Now the columns of the matrix $`[C(\mathrm{\Psi }),I_{sm}]^{\text{tr}}\mathrm{Mat}_{m\times s}(\mathrm{\Lambda })`$ form a basis for $`W`$. Thus,
$$\left[I_mC(\mathrm{\Psi }\gamma ^1)\right]\left[\begin{array}{c}C(\mathrm{\Psi })\\ I_{sm}\end{array}\right]=0,\gamma \mathrm{\Gamma }(E),$$
and so $`C(\mathrm{\Psi }\gamma )=C(\mathrm{\Psi }),\gamma \mathrm{\Gamma }(E)`$. After varying over all $`E/𝔽_q(t)`$ finite, it follows from Theorem 4.4.6 that $`C(\mathrm{\Psi })\mathrm{Mat}_{m\times (sm)}(\overline{k}(t))`$. ∎
###### Proposition 4.5.8.
Let $`M`$ be a $`t`$-motive. Suppose that $`\mathrm{\Phi }\mathrm{GL}_r(\overline{k}(t))`$ represents multiplication by $`𝛔`$ on $`M`$ and that $`\mathrm{\Psi }\mathrm{GL}_r(𝕃)`$ is a rigid analytic trivialization for $`\mathrm{\Phi }`$. For every $`t`$-motive $`N`$ in $`𝒯_M`$ and every $`\mathrm{\Gamma }`$-subrepresentation $`V`$ of $`N^\mathrm{B}`$, there is a sub-$`t`$-motive $`PN`$ so that $`\xi _M(P)=V`$.
###### Proof.
Pick a $`\overline{k}(t)`$-basis $`𝐧\mathrm{Mat}_{s\times 1}`$ for $`N`$ with $`𝝈𝐧=\mathrm{\Phi }_N𝐧`$, and let $`\mathrm{\Psi }_N\mathrm{GL}_s(𝕃)`$ be a rigid analytic trivialization for $`\mathrm{\Phi }_N`$. Let $`𝐯\mathrm{Mat}_{v\times 1}(N^\mathrm{B})`$ be an $`𝔽_q(t)`$-basis for $`V`$, and extend $`𝐯`$ to a basis $`𝐮`$ of $`N^\mathrm{B}`$, $`𝐮=[𝐯,𝐰]^{\text{tr}}`$. By Lemma 4.5.2, there is a $`H(\mathrm{\Psi })\mathrm{GL}_s(\mathrm{\Sigma })`$ so that $`𝐮=H(\mathrm{\Psi })𝐧`$. We note that $`H(\mathrm{\Psi })=\delta ^1\mathrm{\Psi }_N^1`$ for some $`\delta \mathrm{GL}_s(𝔽_q(t))`$ by Proposition 3.3.9(b).
Let $`E/𝔽_q(t)`$ be a finite extension of fields, and let $`\gamma \mathrm{\Gamma }(E)`$. The action of $`\gamma `$ on $`E_{𝔽_q(t)}N^\mathrm{B}`$ is given by the restriction of $`\mathrm{\Xi }^{(E)}`$ as in the proof of Proposition 4.5.6 to $`E_{𝔽_q(t)}N^\mathrm{B}`$. Thus,
$$\xi ^{(E)}(\gamma )(𝐮)=H(\mathrm{\Psi }\gamma )𝐧=H(\mathrm{\Psi }\gamma )H(\mathrm{\Psi })^1𝐮.$$
Since $`V`$ is invariant under $`\mathrm{\Gamma }`$, it follows that the upper right $`v\times (sv)`$ block of $`H(\mathrm{\Psi }\gamma )H(\mathrm{\Psi })^1`$ is $`0`$ for every $`\gamma \mathrm{\Gamma }(E)`$. Let $`D(\mathrm{\Psi })\mathrm{Mat}_{s\times (sv)}(\mathrm{\Lambda })`$ be the $`sv`$ right-most columns of $`H(\mathrm{\Psi })^1`$, and consider the subspace $`W\mathrm{Mat}_{1\times s}(\mathrm{\Lambda })`$,
$$W=\{𝐱\mathrm{Mat}_{1\times s}(\mathrm{\Lambda })𝐱D(\mathrm{\Psi })=0\}.$$
By our considerations on $`H(\mathrm{\Psi })`$ at the end of the preceding paragraph, we see from Lemma 4.5.7 that $`W`$ has a set of defining equations over $`\overline{k}(t)`$. Thus there is a $`C\mathrm{Mat}_{v\times s}(\overline{k}(t))`$ of maximal rank so that $`CD(\mathrm{\Psi })=0`$. Extend $`C`$ to a matrix $`B\mathrm{GL}_s(\overline{k}(t))`$ such that $`C`$ forms the top rows of $`B`$. Now let $`𝐧^{}=B𝐧=[𝐩,𝐪]^{\text{tr}}`$, with $`𝝈𝐧^{}=\mathrm{\Phi }^{}𝐧^{}`$, and let $`P`$ be the $`\overline{k}(t)`$-span of $`𝐩=C𝐧`$. Then
$$𝝈\left[\begin{array}{c}𝐩\\ 𝐪\end{array}\right]=𝝈(B𝐧)=𝝈(BH(\mathrm{\Psi })^1H(\mathrm{\Psi })𝐧)=\left(BH(\mathrm{\Psi })^1\right)^{(1)}H(\mathrm{\Psi })B^1\left[\begin{array}{c}𝐩\\ 𝐪\end{array}\right].$$
By construction, the upper right-hand $`v\times (sv)`$ block of $`BH(\mathrm{\Psi })^1`$ is $`0`$. Thus,
$$𝝈\left[\begin{array}{c}𝐩\\ 𝐪\end{array}\right]=\left[\begin{array}{cc}\mathrm{\Phi }_P& 0\\ & \end{array}\right]\left[\begin{array}{c}𝐩\\ 𝐪\end{array}\right]=\mathrm{\Phi }^{}\left[\begin{array}{c}𝐩\\ 𝐪\end{array}\right].$$
Since $`\mathrm{\Phi }^{}\mathrm{GL}_s(\overline{k}(t))`$, it follows that $`\mathrm{\Phi }_P\mathrm{GL}_v(\overline{k}(t))`$. Thus $`P`$ is a sub-$`t`$-motive of $`N`$. Furthermore, as $`H(\mathrm{\Psi })^1=\mathrm{\Psi }_N\delta `$, $`\delta \mathrm{GL}_s(𝔽_q(t))`$, it follows that $`BH(\mathrm{\Psi })^1`$ is a rigid analytic trivialization of $`\mathrm{\Phi }^{}`$. If we set take $`\mathrm{\Psi }_P`$ to be the upper left-hand block of $`BH(\mathrm{\Psi })^1`$, then $`\mathrm{\Psi }_P`$ is a rigid analytic trivialization for $`\mathrm{\Phi }_P`$. Moreover, it follows that $`P^\mathrm{B}=V`$ by Proposition 3.3.9(b). ∎
###### Proposition 4.5.9.
Let $`M`$ be a $`t`$-motive. Suppose that $`\mathrm{\Phi }\mathrm{GL}_r(\overline{k}(t))`$ represents multiplication by $`𝛔`$ on $`M`$ and that $`\mathrm{\Psi }\mathrm{GL}_r(𝕃)`$ is a rigid analytic trivialization for $`\mathrm{\Phi }`$. To every representation $`W`$ in $`\mathrm{𝐑𝐞𝐩}(\mathrm{\Gamma },𝔽_q(t))`$ there is a $`t`$-motive $`N`$ in $`𝒯_M`$ so that $`W`$ is isomorphic to a subquotient of $`\xi _M(N)`$.
###### Proof.
The representation $`M^\mathrm{B}`$ is faithful by Corollary 4.5.4. Thus any object in $`\mathrm{𝐑𝐞𝐩}(\mathrm{\Gamma },𝔽_q(t))`$ is isomorphic to a subquotient of a direct sum of representations of the form $`(M^\mathrm{B})_v^u:=(M^\mathrm{B})^u((M^\mathrm{B})^{})^v`$. Since $`\xi _M(M_v^u)=(M_v^u)^\mathrm{B}(M^\mathrm{B})_v^u`$, the proposition follows. ∎
###### Theorem 4.5.10.
Let $`M`$ be a $`t`$-motive. Suppose that $`\mathrm{\Phi }\mathrm{GL}_r(\overline{k}(t))`$ represents multiplication by $`𝛔`$ on $`M`$ and that $`\mathrm{\Psi }\mathrm{GL}_r(𝕃)`$ is a rigid analytic trivialization for $`\mathrm{\Phi }`$. Then the functor
$$\xi _M:𝒯_M\mathrm{𝐑𝐞𝐩}(\mathrm{\Gamma },𝔽_q(t))$$
is an equivalence of Tannakian categories. Equivalently, the homomorphism $`\pi _M:\mathrm{\Gamma }\mathrm{\Gamma }_M`$ is an isomorphism over $`𝔽_q(t)`$.
###### Proof.
By Propositions 4.5.6 and 4.5.8, the map $`\pi _M`$ is faithfully flat \[11, Prop. II.2.21(a)\]. By Proposition 4.5.9, $`\pi _M`$ is a closed immersion \[11, Prop. II.2.21(b)\]. Thus $`\pi _M`$ is an isomorphism of affine group schemes over $`𝔽_q(t)`$. ∎
#### 4.5.11. Remark
Although we have focused on objects in the category $`𝒯`$, the above theorem is true (with the same proof) if $`M`$ is replaced by simply a rigid analytically trivial pre-$`t`$-motive.
## 5. Galois groups and transcendence
In this section we first recall the linear independence criterion introduced in by Anderson, Brownawell, and the author, and one of its applications to $`t`$-motives. We then link this together with our study of the Galois groups of certain $`t`$-motives, whose matrices representing multiplication by $`𝝈`$ have entries in $`\overline{k}[t]`$ and whose fundamental matrices have entries in $`𝔼`$. These $`t`$-motives include as a subset rigid analytically trivial Anderson $`t`$-motives. In what follows our primary goal will be to consider the fundamental matrix $`\mathrm{\Psi }`$ associated to such a $`t`$-motive $`M`$ and to equate the transcendence degree over $`\overline{k}`$ of $`\mathrm{\Psi }(\theta )`$ and the dimension of the Galois group of $`M`$.
### 5.1. Linear independence criterion
###### Theorem 5.1.1 (\[2, Thm. 3.1.1\]).
Let $`\mathrm{\Phi }\mathrm{Mat}_r(\overline{k}[t])`$ be given such that $`det\mathrm{\Phi }=c(t\theta )^s`$, $`c\overline{k}^\times `$, and suppose that $`\psi \mathrm{Mat}_{r\times 1}(𝔼)`$ satisfies
$$\psi ^{(1)}=\mathrm{\Phi }\psi .$$
For every $`\rho \mathrm{Mat}_{1\times r}(\overline{k})`$ such that $`\rho \psi (\theta )=0`$, there is a $`P\mathrm{Mat}_{1\times r}(\overline{k}[t])`$ so that $`P(\theta )=\rho `$ and $`P\psi =0`$.
#### 5.1.2. Connection with solutions of $`\sigma `$-semilinear equations
At first glance at the above theorem, the solutions $`\psi `$ of the $`\sigma `$-semilinear equation associated to $`\mathrm{\Phi }`$ are quite special in that their entries are assumed to be in $`𝔼`$. However, the following proposition demonstrates that this situation is not unusual.
###### Proposition 5.1.3 (\[2, Prop. 3.1.3\]).
Suppose we are given $`\mathrm{\Phi }\mathrm{Mat}_r(\overline{k}[t])`$ and $`\psi \mathrm{Mat}_{r\times 1}(𝕋)`$ so that
$$det\mathrm{\Phi }(0)0,\psi ^{(1)}=\mathrm{\Phi }\psi .$$
Then we necessarily have $`\psi \mathrm{Mat}_{r\times 1}(𝔼)`$.
#### 5.1.4. Connection with left $`\overline{k}[t;𝝈]`$-modules
The following is a variation on \[2, Prop. 4.4.3\] with slightly milder hypotheses. We do not assume that the representing matrix $`\mathrm{\Phi }`$ is one directly associated to an Anderson $`t`$-motive. However, we do obtain the same equality of dimensions (with the same proof).
###### Proposition 5.1.5 (\[2, Prop. 4.4.3\]).
Let $`\mathrm{\Phi }\mathrm{Mat}_r(\overline{k}[t])`$ and $`\psi \mathrm{Mat}_{r\times 1}(𝔼)`$ be given as in Theorem 5.1.1. Let $`N`$ be the $`\overline{k}[t]`$-span in $`𝔼`$ of the entries of $`\psi `$, and let $`V`$ be the $`\overline{k}`$-span in $`\overline{k_{\mathrm{}}}`$ of the entries of $`\psi (\theta )`$. Then $`\mathrm{rk}_{\overline{k}[t]}N=dim_{\overline{k}}V`$.
###### Proof.
Let $`N_1:=\{P\mathrm{Mat}_{1\times r}(\overline{k}[t])P\psi =0\}`$. We then obtain an exact sequence of $`\overline{k}[t]`$-modules,
$$0N_1\mathrm{Mat}_{1\times r}(\overline{k}[t])N0,$$
where the second map is given by $`PP\psi `$. It is easy to check that this is an exact sequence of left $`\overline{k}[t;𝝈]`$-modules. Every $`\overline{k}[t]`$-basis for $`N_1`$ can be extended to a basis of $`\mathrm{Mat}_{1\times r}(\overline{k}[t])`$, and so the number of $`\overline{k}`$-linearly independent relations of $`\overline{k}`$-linear dependence among the entries of $`\psi (\theta )`$ is at least as great as $`\mathrm{rk}_{\overline{k}[t]}N_1`$. Thus $`\mathrm{rk}_{\overline{k}[t]}Ndim_{\overline{k}}V`$. Moreover, Theorem 5.1.1 implies that every $`\overline{k}`$-linear relation among the entries of $`\psi (\theta )`$ lifts to a $`\overline{k}[t]`$-linear relation among the entries of $`\psi `$. Thus $`\mathrm{rk}_{\overline{k}[t]}Ndim_{\overline{k}}V`$. ∎
### 5.2. Dimensions and transcendence degrees
#### 5.2.1. Rigid analytic trivializations over $`𝔼`$
Let $`M`$ be a $`t`$-motive. Suppose that $`\mathrm{\Phi }\mathrm{GL}_r(\overline{k}(t))\mathrm{Mat}_r(\overline{k}[t])`$ represents multiplication by $`𝝈`$ on $`M`$ and that $`det\mathrm{\Phi }=c(t\theta )^s`$, $`c\overline{k}^\times `$. An important observation is that by Propositions 3.3.9(c) and 5.1.3, there is a rigid analytic trivialization $`\mathrm{\Psi }`$ for $`\mathrm{\Phi }`$ such that $`\mathrm{\Psi }\mathrm{GL}_r(𝕋)\mathrm{Mat}_r(𝔼)`$.
###### Theorem 5.2.2.
Let $`M`$ be a $`t`$-motive, and let $`\mathrm{\Gamma }_M`$ be its Galois group. Suppose that $`\mathrm{\Phi }\mathrm{GL}_r(\overline{k}(t))\mathrm{Mat}_r(\overline{k}[t])`$ represents multiplication by $`𝛔`$ on $`M`$ and that $`det\mathrm{\Phi }=c(t\theta )^s`$, $`c\overline{k}^\times `$. Let $`\mathrm{\Psi }`$ be a rigid analytic trivialization of $`\mathrm{\Phi }`$ in $`\mathrm{GL}_r(𝕋)\mathrm{Mat}_r(𝔼)`$. Finally, let $`L`$ be the subfield of $`\overline{k_{\mathrm{}}}`$ generated over $`\overline{k}`$ by the entries of $`\mathrm{\Psi }(\theta )`$. Then
$$\text{tr. deg}_{\overline{k}}L=dim\mathrm{\Gamma }_M.$$
###### Proof.
By Theorem 4.5.10, the groups $`\mathrm{\Gamma }_M`$ and $`\mathrm{\Gamma }_\mathrm{\Psi }`$ are isomorphic. Moreover, by Theorem 4.3.1, their dimension is the same as $`\text{tr. deg}_{\overline{k}(t)}\mathrm{\Lambda }`$, where $`\mathrm{\Lambda }=\overline{k}(t)(\mathrm{\Psi })𝕃`$. Now let $`Q=\overline{k}[\mathrm{\Psi }(\theta )]L`$, and let $`S=\overline{k}(t)[\mathrm{\Psi }]\mathrm{\Lambda }`$. Then as rings,
$$Q\overline{k}[X_{ij}]/𝔞,S\overline{k}(t)[X_{ij}]/𝔟,$$
for ideals $`𝔞`$ and $`𝔟`$. For $`d1`$, let $`\overline{k}[X_{ij}]_d`$ and $`𝔞_d`$ denote the elements of $`\overline{k}[X_{ij}]`$ and $`𝔞`$ of total degree $`d`$, and let $`Q_dQ`$ correspond to their quotient. Similarly define $`\overline{k}(t)[X_{ij}]_d`$, $`𝔟_d`$, and $`S_d`$.
Fix $`d1`$. Now for any $`n1`$, the entries of $`\mathrm{\Psi }^n`$ comprise all monomials of total degree $`n`$ in the $`\mathrm{\Psi }_{ij}`$. If $`\psi `$ is a column of $`\mathrm{\Psi }^n`$, then $`\psi ^{(1)}=\mathrm{\Phi }^n\psi `$. Thus let $`\overline{\psi }\mathrm{Mat}_{N\times 1}(𝔼)`$ be the column vector whose entries are the concatenation of $`1`$ and each of the columns of $`\mathrm{\Psi }^n`$ for $`nd`$. (Here $`N=(r^{2d+2}1)/(r^21)`$.) Then if $`\overline{\mathrm{\Phi }}\mathrm{Mat}_N(\overline{k}[t])\mathrm{GL}_N(\overline{k}(t))`$ is the block diagonal matrix
$$\overline{\mathrm{\Phi }}:=[1]\mathrm{\Phi }^r\left(\mathrm{\Phi }^2\right)^{r^2}\mathrm{}\left(\mathrm{\Phi }^d\right)^{r^d},$$
it follows that $`\overline{\psi }^{(1)}=\overline{\mathrm{\Phi }}\overline{\psi }`$. Now it is easy to see that $`Q_d`$ is the $`\overline{k}`$-span of the columns of $`\overline{\psi }(\theta )`$ and that $`S_d`$ is the $`\overline{k}(t)`$-span of the columns of $`\overline{\psi }`$. Since $`\overline{\mathrm{\Phi }}`$ and $`\overline{\psi }`$ satisfy the hypotheses for Proposition 5.1.5, we see that for all $`d1`$,
$$dim_{\overline{k}}Q_d=dim_{\overline{k}(t)}S_d.$$
Thus the homogenizations of $`Q`$ and $`S`$ have the same Hilbert series (see \[34, Ch. VII, §12\]), and so $`\text{tr. deg}_{\overline{k}}L=\text{tr. deg}_{\overline{k}(t)}\mathrm{\Lambda }`$. ∎
## 6. Application to Carlitz logarithms
### 6.1. Carlitz logarithms and $`t`$-motives
#### 6.1.1. The power series $`L_\alpha `$
For $`\alpha \overline{k}^\times `$ with $`\left|\alpha \right|_{\mathrm{}}<\left|\theta \right|_{\mathrm{}}^{q/(q1)}`$, define the power series
$$L_\alpha (t):=\alpha +\underset{i=1}{\overset{\mathrm{}}{}}\frac{\alpha ^{q^i}}{(t\theta ^q)(t\theta ^{q^2})\mathrm{}(t\theta ^{q^i})}.$$
It is easy to show that $`L_\alpha 𝕋`$ and that moreover, $`L_\alpha (z)`$ converges for all $`z𝕂`$ with $`\left|z\right|_{\mathrm{}}<\left|\theta \right|_{\mathrm{}}^q`$. By §1.2.2, we see that
$$L_\alpha (\theta )=\mathrm{log}_C(\alpha ).$$
Furthermore, as a power series in $`𝕋`$, $`L_\alpha `$ also satisfies the functional equation
(6.1.1.1)
$$L_\alpha ^{(1)}=\alpha ^{(1)}+\frac{L_\alpha }{t\theta }.$$
#### 6.1.2. $`t`$-motives for Carlitz logarithms
Fix $`\alpha _1,\mathrm{},\alpha _r\overline{k}^\times `$ with $`\left|\alpha _i\right|_{\mathrm{}}<\left|\theta \right|_{\mathrm{}}^{q/(q1)}`$ for $`i=1,\mathrm{},r`$. Set
$$\mathrm{\Phi }:=\mathrm{\Phi }(\alpha _1,\mathrm{},\alpha _r):=\left[\begin{array}{cccc}t\theta & 0& \mathrm{}& 0\\ \alpha _1^{(1)}(t\theta )& 1& \mathrm{}& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ \alpha _r^{(1)}(t\theta )& 0& \mathrm{}& 1\end{array}\right]\mathrm{Mat}_{r+1}(\overline{k}[t]).$$
Note that $`\mathrm{\Phi }`$ defines a pre-$`t`$-motive $`X:=X(\alpha _1,\mathrm{},\alpha _r)`$ that is an extension of $`\mathrm{𝟏}^r`$ by the Carlitz motive $`C`$:
$$0CX\mathrm{𝟏}^r0.$$
In spite of the restrictions on $`\alpha _1,\mathrm{},\alpha _r`$, we will be able to use the objects $`X(\alpha _1,\mathrm{},\alpha _r)`$ to accommodate *all* Carlitz logarithms using Lemma 6.4.1.
###### Proposition 6.1.3.
Let $`\alpha _1,\mathrm{},\alpha _r\overline{k}^\times `$ with $`\left|\alpha _i\right|_{\mathrm{}}<\left|\theta \right|_{\mathrm{}}^{q/(q1)}`$ for $`i=1,\mathrm{},r`$.The pre-$`t`$-motive $`X=X(\alpha _1,\mathrm{},\alpha _r)`$ is a $`t`$-motive.
###### Proof.
We prove first that $`X`$ is rigid analytically trivial and then that $`X`$ is an object in $`𝒯`$. Define
$$\mathrm{\Psi }:=\mathrm{\Psi }(\alpha _1,\mathrm{},\alpha _r):=\left[\begin{array}{cccc}\mathrm{\Omega }& 0& \mathrm{}& 0\\ \mathrm{\Omega }L_{\alpha _1}& 1& \mathrm{}& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ \mathrm{\Omega }L_{\alpha _r}& 0& \mathrm{}& 1\end{array}\right]\mathrm{GL}_{r+1}(𝕋)$$
It is a simple matter to check that $`\mathrm{\Psi }`$ is a rigid analytic trivialization for $`\mathrm{\Phi }`$ using (6.1.1.1). We note by Proposition 5.1.3 that the entries of $`\mathrm{\Psi }`$ are in $`𝔼`$.
Consider the pre-$`t`$-motive $`CX`$. We claim that $`CX`$ is in the essential image of the functor $`𝖬M:𝒜^I`$ of Theorem 3.4.9. By the definition of the category $`𝒯`$ in §3.4.10, it will follow that $`X`$ is a $`t`$-motive.
Let $`𝖬:=\overline{k}[t]^{r+1}`$ with standard $`\overline{k}[t]`$-basis $`m_0,\mathrm{},m_r`$. Letting $`𝗆:=[m_1,\mathrm{},m_r]^{\text{tr}}`$, we give $`𝖬`$ the structure of a left $`\overline{k}[t;𝝈]`$-module by setting
$$𝝈𝗆:=(t\theta )\mathrm{\Phi }𝗆.$$
Now $`𝖬`$ sits in an exact sequence of left $`\overline{k}[t;𝝈]`$-modules,
$$0𝖢^2𝖬𝖢^r0,$$
where $`𝖢`$ is the Carlitz motive in the category of Anderson $`t`$-motives of §3.4.3. Since $`𝖢`$ and $`𝖢^{}`$ are finitely generated as left $`\overline{k}[\sigma ]`$-modules, so is $`𝖬`$, and it follows from \[2, Prop. 4.3.2\] that $`𝖬`$ is free and finitely generated as a left $`\overline{k}[𝝈]`$-module. Without much difficulty, one shows that $`𝝈𝖬=(t\theta )^2m_0,(t\theta )m_1,\mathrm{},(t\theta )m_r_{\overline{k}[t]}`$. Thus $`(t\theta )^n𝖬𝝈𝖬`$ for all $`n2`$, and $`𝖬`$ is an Anderson $`t`$-motive by §3.4.1. ∎
### 6.2. The Galois group $`\mathrm{\Gamma }_X`$
We continue with the notations of the previous section, including choices of $`\alpha _1,\mathrm{},\alpha _r\overline{k}^\times `$ with $`\left|\alpha _i\right|_{\mathrm{}}<\left|\theta \right|_{\mathrm{}}^{q/(q1)}`$ for $`i=1,\mathrm{},r`$.
#### 6.2.1. The group $`G`$
Let $`G`$ be the algebraic subgroup of $`\mathrm{GL}_{r+1}`$ over $`𝔽_q(t)`$ such that for all $`𝔽_q(t)`$-algebras $`R`$,
$$G(R)=\left\{\left[\begin{array}{cc}& 0\\ & I_r\end{array}\right]\mathrm{GL}_{r+1}(R)\right\}.$$
#### 6.2.2. Preliminary calculations
We claim that $`\mathrm{\Gamma }_XG`$. As in §4.2, we can construct the coordinate ring as the image of $`\mu :𝔽_q(t)[X,1/detX]𝕃_{\overline{k}(t)}𝕃`$, the $`𝔽_q(t)`$-algebra homomorphism that sends $`X`$ to $`\stackrel{~}{\mathrm{\Psi }}=\mathrm{\Psi }_1^1\mathrm{\Psi }_2`$. As before let $`𝔮=\mathrm{ker}\mu `$. Direct calculation verifies that $`X_{ij}\delta _{ij}𝔮`$ for all $`i1`$ and $`j2`$, where $`\delta _{ij}`$ is the usual Kronecker delta. Thus, $`\mathrm{\Gamma }_XG`$. It will be convenient henceforth to label the non-trivial coordinates of $`G\mathrm{GL}_{r+1}`$ as $`X_0,\mathrm{},X_r`$.
Because the Carlitz motive $`C`$ is contained in $`X`$, it is an object in $`𝒯_X`$, and hence there is a surjection $`\pi :\mathrm{\Gamma }_X𝔾_\mathrm{m}`$ over $`𝔽_q(t)`$ by Theorem 3.5.4. Now under $`\nu :\overline{k}(t)[X_0,X_0^1,X_1,\mathrm{},X_r]𝕃`$, which takes $`X`$ to $`\mathrm{\Psi }`$, we have $`\nu (X_0)=\mathrm{\Omega }`$. Thus the action of any $`\gamma \mathrm{\Gamma }_X(\overline{𝔽_q(t)})`$ on $`\mathrm{\Omega }`$ agrees with the action of the $`X_0`$-coordinate of $`\gamma `$ on $`\mathrm{\Omega }`$. That is, the surjection $`\pi `$ coincides with the natural projection on the $`X_0`$-coordinate of $`G`$. Let $`V`$ be the kernel of $`\pi `$ so that we have an exact sequence of algebraic groups over $`𝔽_q(t)`$,
$$1V\mathrm{\Gamma }_X𝔾_\mathrm{m}1.$$
The group $`V`$ is a subgroup of the group of unipotent matrices of $`G`$, which itself is naturally isomorphic to $`𝔾_\mathrm{a}^r`$. Thus we can think of $`V𝔾_\mathrm{a}^r`$ with coordinates $`X_1,\mathrm{},X_r`$.
###### Proposition 6.2.3.
With notation as above, the group $`V`$ is a linear subspace of $`𝔾_\mathrm{a}^r`$ over $`𝔽_q(t)`$.
###### Proof.
Since $`\mathrm{\Gamma }_X`$ is a smooth over $`𝔽_q(t)`$ by Theorem 4.3.1, one verifies that the map $`\pi `$ is surjective on Lie algebras, and hence $`V`$ is also smooth. Thus it is determined by the Zariski closure of $`V(\overline{𝔽_q(t)})`$ in $`𝔾_\mathrm{a}^r`$. Because $`\pi `$ is surjective, for any non-zero $`\alpha \overline{𝔽_q(t)}`$, we can choose $`\gamma \mathrm{\Gamma }_X(\overline{𝔽_q(t)})`$ so that $`\pi (\gamma )=\alpha `$. Suppose that $`\mu =\left[\begin{array}{cc}1& 0\\ v& I_r\end{array}\right]V(\overline{𝔽_q(t)})`$. Then direct calculation gives $`\gamma ^1\mu \gamma =\left[\begin{array}{cc}1& 0\\ \alpha v& I_r\end{array}\right]V(\overline{𝔽_q(t)})`$, and thus $`V(\overline{𝔽_q(t)})`$ is a linear subspace of $`𝔾_\mathrm{a}^r(\overline{𝔽_q(t)})`$. Since $`V`$ is smooth, its defining equations over $`\overline{𝔽_q(t)}`$ are linear forms in $`X_1,\mathrm{},X_r`$. These forms can be defined over $`𝔽_q(t)`$ since $`V`$ is simply a linear subspace. ∎
#### 6.2.4. Defining polynomials for $`\mathrm{\Gamma }_X`$
Because the map $`\mathrm{\Gamma }_X𝔾_\mathrm{m}`$ is a smooth morphism over $`𝔽_q(t)`$, Hilbert’s Theorem 90 provides an exact sequence
$$1V(𝔽_q(t))\mathrm{\Gamma }_X(𝔽_q(t))𝔾_\mathrm{m}(𝔽_q(t))1$$
by \[32, §18.5\]. Let $`b_0𝔽_q(t)^\times 𝔽_q^\times `$, and fix a matrix
(6.2.4.1)
$$\gamma =\left[\begin{array}{cccc}b_0& 0& \mathrm{}& 0\\ b_1& 1& \mathrm{}& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ b_r& 0& \mathrm{}& 1\end{array}\right]\mathrm{\Gamma }_X(𝔽_q(t))$$
One checks that the Zariski closure in $`\mathrm{\Gamma }_X`$ of the cyclic group generated by $`\gamma `$ is the line in $`G`$ connecting $`\gamma `$ to the identity matrix. Translating this line by any element of $`V`$ shows that $`\mathrm{\Gamma }_X`$ contains the linear space spanned by $`V`$ and $`\gamma `$. Since $`\mathrm{\Gamma }_X`$ is irreducible and of dimension $`1`$ greater than the dimension of $`V`$, we see that $`\mathrm{\Gamma }_X`$ is this linear subspace. Moreover, this implies the following proposition.
###### Proposition 6.2.5.
Suppose $`F_1,\mathrm{},F_s𝔽_q(t)[X_1,\mathrm{},X_r]`$ are linear forms defining $`V`$, and suppose that $`\gamma \mathrm{\Gamma }_X(𝔽_q(t))`$ is defined as in (6.2.4.1). Then the linear polynomials in $`𝔽_q(t)[X_0,\mathrm{},X_r]`$,
$$G_i:=(b_01)F_iF_i(b_1,\mathrm{},b_r)(X_01),i=1,\mathrm{},s,$$
are defining polynomials for $`\mathrm{\Gamma }_X`$.
### 6.3. Linear relations among Carlitz logarithms
#### 6.3.1. Defining polynomials for $`Z`$
As usual let $`Z:=\mathrm{Spec}\mathrm{\Sigma }`$, where $`\mathrm{\Psi }=\mathrm{\Psi }(\alpha _1,\mathrm{},\alpha _m)`$. From Proposition 4.2.8 we see that $`Z`$ and $`\mathrm{\Gamma }_X`$ are isomorphic over $`\overline{k(t)}`$. Since $`\mathrm{\Gamma }_X`$ is a linear space, $`Z`$ is also a linear space and isomorphic to $`\mathrm{\Gamma }_X`$ over $`\overline{k}(t)`$. Thus we can pick
$$\zeta =\left[\begin{array}{cccc}f_0& 0& \mathrm{}& 0\\ f_1& 1& \mathrm{}& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ f_r& 0& \mathrm{}& 1\end{array}\right]Z(\overline{k}(t)),$$
and then
$$Z(\overline{k}(t))=\zeta \mathrm{\Gamma }_X(\overline{k}(t)).$$
It is a simple matter to check that the linear polynomials in $`\overline{k}(t)[X_0,\mathrm{},X_r]`$,
$$H_i:=G_iX_0G(f_0,\mathrm{},f_r)/f_0,i=1,\mathrm{},s,$$
are defining polynomials for $`Z`$.
The following theorems show how the above constructions can be used to characterize all $`k`$-linear relations among $`\stackrel{~}{\pi }`$, $`\mathrm{log}_C(\alpha _1),\mathrm{},\mathrm{log}_C(\alpha _r)`$.
###### Theorem 6.3.2.
Let $`\alpha _1,\mathrm{},\alpha _r\overline{k}^\times `$ with $`\left|\alpha _i\right|_{\mathrm{}}<\left|\theta \right|_{\mathrm{}}^{q/(q1)}`$ for $`i=1,\mathrm{},r`$. Let $`X=X(\alpha _1,\mathrm{},\alpha _r)`$ be the associated $`t`$-motive.
1. Let $`F=c_1X_1+\mathrm{}+c_rX_r`$, $`c_1,\mathrm{},c_r𝔽_q(t)`$, be a defining linear form for $`V`$ so that $`G=(b_01)FF(b_1,\mathrm{},b_r)(X_01)`$, $`b_0,\mathrm{},b_r𝔽_q(t),b_0𝔽_q`$, is a defining polynomial for $`\mathrm{\Gamma }_X`$. Then
$$(b_0(\theta )1)\underset{i=1}{\overset{r}{}}c_i(\theta )\mathrm{log}_C(\alpha _i)\underset{i=1}{\overset{r}{}}c_i(\theta )b_i(\theta )\stackrel{~}{\pi }=0.$$
2. Every $`k`$-linear relation among $`\stackrel{~}{\pi }`$, $`\mathrm{log}_C(\alpha _1),\mathrm{},\mathrm{log}_C(\alpha _r)`$ is a $`k`$-linear combination of the relations from part (a).
3. Let $`N`$ be the $`k`$-linear span of $`\stackrel{~}{\pi }`$, $`\mathrm{log}_C(\alpha _1),\mathrm{},\mathrm{log}_C(\alpha _r)`$. Then $`dim\mathrm{\Gamma }_X=dim_kN`$.
###### Proof.
Choose $`f\overline{k}(t)`$ as in §6.3.1 so that $`H:=GfX_0`$ is a defining polynomial for $`Z`$. Then
(6.3.2.1)
$$H(\mathrm{\Omega },\mathrm{\Omega }L_{\alpha _1},\mathrm{},\mathrm{\Omega }L_{\alpha _r})=G(\mathrm{\Omega },\mathrm{\Omega }L_{\alpha _1},\mathrm{},\mathrm{\Omega }L_{\alpha _r})f\mathrm{\Omega }=0.$$
We see that
$$\begin{array}{c}f^{(1)}\mathrm{\Omega }^{(1)}=\sigma G(\mathrm{\Omega },\mathrm{\Omega }L_{\alpha _1},\mathrm{},\mathrm{\Omega }L_{\alpha _r})=\mathrm{\Omega }G(t\theta 1,\alpha _1^{(1)}(t\theta ),\mathrm{},\alpha _r^{(1)}(t\theta ))\hfill \\ \hfill +f\mathrm{\Omega }F(b_1,\mathrm{},b_r)\mathrm{\Omega }.\end{array}$$
The first equality is a consequence of (6.3.2.1), and the second follows from direct computation. Thus
$$(t\theta )f^{(1)}f=G(t\theta 1,\alpha _1^{(1)}(t\theta ),\mathrm{},\alpha _r^{(1)}(t\theta ))F(b_1,\mathrm{},b_r).$$
The right-hand side is a polynomial in $`\overline{k}[t]`$, so it follows that $`f`$ is regular at $`t=\theta `$. Indeed if not, then $`f^{(1)}`$ must have a pole at $`t=\theta ^{(1)}`$, whence $`f`$ must also have a pole at $`t=\theta ^{(1)}`$. Continuing in this way we see that if $`f`$ has a pole at $`t=\theta `$, then it must have a pole at each $`t=\theta ^{(i)}`$, $`i1`$, which is not possible. By a similar argument we deduce that $`f^{(1)}`$ is also regular at $`t=\theta `$. Thus we see that
$$f(\theta )=G(1,0,\mathrm{},0)|_{t=\theta }+\underset{i=1}{\overset{r}{}}c_i(\theta )b_i(\theta )=\underset{i=1}{\overset{r}{}}c_i(\theta )b_i(\theta ).$$
Equation (6.3.2.1) transforms into
$$(b_01)\underset{i=1}{\overset{r}{}}c_i\mathrm{\Omega }L_{\alpha _i}\underset{i=1}{\overset{r}{}}c_ib_i(\mathrm{\Omega }1)f\mathrm{\Omega }=0.$$
Dividing through by $`\mathrm{\Omega }`$ and evaluating at $`t=\theta `$, we obtain part (a). Part (b) is a consequence of (a) and (c), since $`\mathrm{\Gamma }_X`$ is a linear space in $`G`$ over $`𝔽_q(t)`$. For part (c), part (a) implies that $`dim_kNdim\mathrm{\Gamma }_X`$, since the defining polynomials for $`\mathrm{\Gamma }_X`$ generate a set of $`k`$-linear relations on $`\stackrel{~}{\pi }`$, $`\mathrm{log}_C(\alpha _1),\mathrm{},\mathrm{log}_C(\alpha _r)`$ of dimension $`r+1dim\mathrm{\Gamma }_X`$. However, $`dim_kN\text{tr. deg}_{\overline{k}}\overline{k}(\stackrel{~}{\pi },\mathrm{log}_C(\alpha _1),\mathrm{},\mathrm{log}_C(\alpha _r))`$ and the latter quantity is $`dim\mathrm{\Gamma }_X`$ by Theorem 5.2.2. ∎
#### 6.3.3. Example
Let $`\zeta _\theta =\sqrt[q1]{\theta }`$, let $`X`$ be the $`t`$-motive $`X(\zeta _\theta )`$ of dimension $`2`$ over $`\overline{k}(t)`$, and let $`\mathrm{\Psi }=\mathrm{\Psi }(\zeta _\theta )`$. Since $`\zeta _\theta `$ satisfies $`_t(\zeta _\theta )=\theta \zeta _\theta +\zeta _\theta ^q=0`$, we see that $`\zeta _\theta `$ is a $`t`$-torsion point on the Carlitz module. Moreover, $`\mathrm{exp}_C(\theta \mathrm{log}_C(\zeta _\theta ))=0`$, and one calculates that
$$\mathrm{log}_C(\zeta _\theta )=\frac{\stackrel{~}{\pi }}{\theta }.$$
Thus $`\mathrm{\Gamma }_X`$ is $`1`$-dimensional by Theorem 6.3.2(c). If we consider the function in $`𝕋`$
$$\mathrm{{\rm Y}}:=tL_{\zeta _\theta }\zeta _\theta (t\theta ),$$
then $`\mathrm{{\rm Y}}^{(1)}=\mathrm{{\rm Y}}/(t\theta )`$. Thus $`\mathrm{{\rm Y}}=f/\mathrm{\Omega }`$ for some $`f𝔽_q[t]`$ by Lemma 3.3.5. Evaluation at $`t=\theta `$ shows that $`f=1`$ identically. Therefore, $`Z_\mathrm{\Psi }`$ is defined by
$$Z_\mathrm{\Psi }:\zeta _\theta (t\theta )X_0tX_11=0.$$
It follows that the defining equation for $`\mathrm{\Gamma }_X`$ is
$$\mathrm{\Gamma }_X:tX_1X_0+1=0.$$
In the notation of Theorem 6.3.2, we have
$$F:=X_1,b_0:=t+1,b_1:=1$$
$$G:=tX_1X_0+1,H:=GfX_0,f:=\zeta _\theta (t\theta )1.$$
### 6.4. Algebraic independence of Carlitz logarithms
Before proving the main result on Carlitz logarithms, we prove a reduction lemma.
###### Lemma 6.4.1.
Let $`\lambda 𝕂^\times `$. If $`\mathrm{exp}_C(\lambda )\overline{k}^\times `$, then there is an $`\alpha \overline{k}^\times `$ with $`\left|\alpha \right|_{\mathrm{}}<\left|\theta \right|_{\mathrm{}}^{q/(q1)}`$, an $`f𝔽_q[\theta ]`$, and an $`n1`$, so that $`\lambda =\theta ^n\mathrm{log}_C(\alpha )+f\stackrel{~}{\pi }`$.
###### Proof.
Let $`\beta =\mathrm{exp}_C(\lambda )`$, and assume that $`\left|\beta \right|_{\mathrm{}}\left|\theta \right|_{\mathrm{}}^{q/(q1)}`$. We solve the equation $`_t(x)=\theta x+x^q=\beta `$; that is, we find the $`t`$-division points of $`\beta `$ on the Carlitz module. The Newton polygon for this equation, along with our assumptions on $`\beta `$, imply that any solution $`\alpha \overline{k}^\times `$ of this equation must satisfy $`\left|\alpha \right|_{\mathrm{}}=\left|\beta \right|_{\mathrm{}}^{1/q}`$. Moreover, if for some $`\eta 𝕂`$ we have $`\mathrm{exp}_C(\eta )=\alpha `$, then
$$\mathrm{exp}_C(\theta \eta )=\beta =\mathrm{exp}_C(\lambda ).$$
If $`\left|\beta \right|_{\mathrm{}}<\left|\theta \right|_{\mathrm{}}^{q^2/(q1)}`$, then $`\alpha `$ is sufficiently small and we can pick $`\eta =\mathrm{log}_C(\alpha )`$. The result then follows with $`n=1`$. Otherwise, we continue to take $`t`$-division values, and for some $`n1`$, we have $`_{t^n}(\alpha )=\beta `$ with $`\left|\alpha \right|_{\mathrm{}}<\left|\theta \right|_{\mathrm{}}^{q/(q1)}`$, for which $`\mathrm{exp}_C(\theta ^n\mathrm{log}_C(\alpha ))=\beta `$. ∎
###### Theorem 6.4.2.
Let $`\lambda _1,\mathrm{},\lambda _r𝕂`$ satisfy $`\mathrm{exp}_C(\lambda _i)\overline{k}`$ for $`i=1,\mathrm{},r`$. If $`\lambda _1,\mathrm{},\lambda _r`$ are linearly independent over $`k`$, then they are algebraically independent over $`\overline{k}`$.
###### Proof.
Assume that $`\lambda _1,\mathrm{},\lambda _r`$ are linearly independent over $`k`$. By Lemma 6.4.1, for each $`\lambda _i`$ we can pick $`\alpha _i\overline{k}^\times `$ with $`\left|\alpha _i\right|_{\mathrm{}}<\left|\theta \right|_{\mathrm{}}^{q/(q1)}`$ so that the $`k`$-linear span of $`\lambda _1,\mathrm{},\lambda _r`$ is contained in the $`k`$-linear span of $`\stackrel{~}{\pi }`$, $`\mathrm{log}_C(\alpha _1),\mathrm{},\mathrm{log}_C(\alpha _r)`$. Let $`X=X(\alpha _1,\mathrm{},\alpha _r)`$ be the $`t`$-motive associated to these logarithms as in the previous sections, and let $`\mathrm{\Gamma }_X`$ be its Galois group. Let
$$L=\overline{k}(\stackrel{~}{\pi },\mathrm{log}_C(\alpha _1),\mathrm{},\mathrm{log}_C(\alpha _r)),$$
and let
$$N=k\text{-linear span of }\stackrel{~}{\pi }\text{}\mathrm{log}_C(\alpha _1),\mathrm{},\mathrm{log}_C(\alpha _r).$$
Because $`\lambda _1,\mathrm{},\lambda _r`$ are linearly independent over $`k`$, we see that $`rdim_kNr+1`$. Theorems 5.2.2 and 6.3.2 imply that
$$\text{tr. deg}_{\overline{k}}L=dim\mathrm{\Gamma }_X=dim_kN.$$
If $`\stackrel{~}{\pi }`$, $`\mathrm{log}_C(\alpha _1),\mathrm{},\mathrm{log}_C(\alpha _r)`$ are linearly independent over $`k`$, then they are algebraically independent over $`\overline{k}`$, whence the same follows for $`\lambda _1,\mathrm{},\lambda _r`$ since $`L=\overline{k}(\stackrel{~}{\pi },\lambda _1,\mathrm{},\lambda _r)`$. If there is a linear dependence among $`\stackrel{~}{\pi }`$, $`\mathrm{log}_C(\alpha _1),\mathrm{},\mathrm{log}_C(\alpha _r)`$ over $`k`$, then $`N`$ is equal to the $`k`$-span of $`\lambda _1,\mathrm{},\lambda _r`$ and $`L=\overline{k}(\lambda _1,\mathrm{},\lambda _r)`$. Thus in that case $`\lambda _1,\mathrm{},\lambda _r`$ are algebraically independent over $`\overline{k}`$. ∎ |
warning/0506/physics0506129.html | ar5iv | text | # 1 Introduction
## 1 Introduction
The process of homogenization involves the combination of two (or more) component material phases to produce a single, effectively homogeneous, composite medium . Typically, the constitutive properties of the component material phases are relatively simple as compared with those of the homogenized composite medium (HCM). Through homogenization, novel and potentially useful material properties may be realized . Many examples of material properties being extended — or indeed entirely new material properties being realized — as a result of homogenization can be found within the regimes of linear and nonlinear electromagnetics .
An interesting result concerns the electromagnetic group velocity in HCMs. Under certain circumstances, the group velocity in an HCM can exceed the group velocities in its component material phases. This issue has been investigated for isotropic dielectric composite mediums using the Maxwell Garnett and the Bruggeman homogenization formalisms. In these studies, an enhancement in group velocity is demonstrated through homogenizing two component material phases, one of which is characterized by a relatively large permittivity and relatively small frequency–dispersive term as compared with the other component material phase.
Enhancement of group velocity in a laminate composite medium has been considered by using a volume–weighted sum to estimate the HCM permittivity . The directional properties of group–velocity enhancement are further explored in this communication. Specifically, we consider a uniaxial dielectric HCM which develops from the homogenization of a random assembly of oriented spheroidal particles. The component material phases are themselves electromagnetically isotropic. Our theoretical analysis is founded upon the Bruggeman homogenization formalism .
## 2 Homogenization
Let us consider the homogenization of a composite medium containing two component material phases, labelled as $`a`$ and $`b`$. Both component material phases are taken to be isotropic dielectric mediums: $`ϵ^a`$ and $`ϵ^b`$ denote the permittivity scalars of phases $`a`$ and $`b`$, respectively. In order to focus in particular upon the phenomenon of enhancement of group velocity, without being distracted by the complications arising from dielectric loss, the component material phases are assumed to be nondissipative; i.e., $`ϵ^{a,b}`$. The component material phases are envisioned as random distributions of identically oriented, spheroidal particles. The spheroidal shape — which is taken to be the same for all particles of phases $`a`$ and $`b`$ — is parameterized via the shape dyadic
$$\underset{¯}{\underset{¯}{U}}=U_{}\underset{¯}{\underset{¯}{I}}+\left(U_{}U_{}\right)\underset{¯}{\overset{^}{c}}\underset{¯}{\overset{^}{c}},$$
(1)
where $`\underset{¯}{\underset{¯}{I}}`$ is the identity 3$`\times `$3 dyadic and the unit vector $`\underset{¯}{\overset{^}{c}}`$ is parallel to the spheroid’s axis of rotational symmetry. The spheroid’s surface is described by the vector
$$\underset{¯}{r}_s(\theta ,\varphi )=\eta \underset{¯}{\underset{¯}{U}}\text{ }\text{}\text{ }\underset{¯}{\overset{^}{r}}(\theta ,\varphi ),$$
(2)
with $`\underset{¯}{\overset{^}{r}}`$ being the radial unit vector from the spheroid’s centroid and specified by the spherical polar coordinates $`\theta `$ and $`\varphi `$. The linear dimensions of the spheroid, as determined by the parameter $`\eta `$, are assumed to be small relative to the electromagnetic wavelength(s).
The permittivity dyadic of the resulting HCM,
$$\underset{¯}{\underset{¯}{ϵ}}^{Br}=ϵ_{}^{Br}\underset{¯}{\underset{¯}{I}}+\left(ϵ_{}^{Br}ϵ_{}^{Br}\right)\underset{¯}{\overset{^}{c}}\underset{¯}{\overset{^}{c}},$$
(3)
is estimated using the Bruggeman homogenization formalism as the solution of the equation
$$f_a\underset{¯}{\underset{¯}{a}}^a+f_b\underset{¯}{\underset{¯}{a}}^b=\underset{¯}{\underset{¯}{0}},$$
(4)
where $`f_a`$ and $`f_b=1f_a`$ denote the respective volume fractions of the material component phases $`a`$ and $`b`$ . The polarizability dyadics in (4) are defined as
$$\underset{¯}{\underset{¯}{a}}^{\mathrm{}}=\left(ϵ^{\mathrm{}}\underset{¯}{\underset{¯}{I}}\underset{¯}{\underset{¯}{ϵ}}^{Br}\right)\text{ }\text{}\text{ }\left[\underset{¯}{\underset{¯}{I}}+\underset{¯}{\underset{¯}{D}}\text{ }\text{}\text{ }\left(ϵ^{\mathrm{}}\underset{¯}{\underset{¯}{I}}\underset{¯}{\underset{¯}{ϵ}}^{Br}\right)\right]^1,(\mathrm{}=a,b),$$
(5)
wherein the depolarization dyadic is given by the surface integral
$$\underset{¯}{\underset{¯}{D}}=\frac{1}{4\pi }_0^{2\pi }𝑑\varphi _0^\pi 𝑑\theta \mathrm{sin}\theta \left(\frac{1}{\underset{¯}{\overset{^}{r}}\text{ }\text{}\text{ }\underset{¯}{\underset{¯}{U}}^1\text{ }\text{}\text{ }\underset{¯}{\underset{¯}{ϵ}}^{Br}\text{ }\text{}\text{ }\underset{¯}{\underset{¯}{U}}^1\text{ }\text{}\text{ }\underset{¯}{\overset{^}{r}}}\right)\underset{¯}{\underset{¯}{U}}^1\text{ }\text{}\text{ }\underset{¯}{\overset{^}{r}}\underset{¯}{\overset{^}{r}}\text{ }\text{}\text{ }\underset{¯}{\underset{¯}{U}}^1.$$
(6)
The depolarization dyadic may be expressed as
$$\underset{¯}{\underset{¯}{D}}=D_{}\underset{¯}{\underset{¯}{I}}+\left(D_{}D_{}\right)\underset{¯}{\overset{^}{c}}\underset{¯}{\overset{^}{c}},$$
(7)
where
$`D_{}`$ $`=`$ $`{\displaystyle \frac{\gamma }{ϵ_{}^{Br}}}\mathrm{\Gamma }_{}(\gamma ),`$ (8)
$`D_{}`$ $`=`$ $`{\displaystyle \frac{1}{ϵ_{}^{Br}}}\mathrm{\Gamma }_{}(\gamma ),`$ (9)
The terms $`\mathrm{\Gamma }_{}`$ and $`\mathrm{\Gamma }_{}`$ herein are functions of the real–valued parameter
$$\gamma =\frac{U_{}^2ϵ_{}^{Br}}{U_{}^2ϵ_{}^{Br}};$$
(10)
they have the representations
$`\mathrm{\Gamma }_{}(\gamma )`$ $`=`$ $`{\displaystyle \frac{1}{4\pi }}{\displaystyle _0^{2\pi }}𝑑\varphi {\displaystyle _0^\pi }𝑑\theta {\displaystyle \frac{\mathrm{cos}^2\varphi \mathrm{sin}^3\theta }{\mathrm{cos}^2\theta +\mathrm{sin}^2\theta \left(\gamma \mathrm{cos}^2\varphi +\mathrm{sin}^2\varphi \right)}},`$ (11)
$`\mathrm{\Gamma }_{}(\gamma )`$ $`=`$ $`{\displaystyle \frac{1}{4\pi }}{\displaystyle _0^{2\pi }}𝑑\varphi {\displaystyle _0^\pi }𝑑\theta {\displaystyle \frac{\mathrm{sin}^2\varphi \mathrm{sin}^3\theta }{\mathrm{cos}^2\theta +\mathrm{sin}^2\theta \left(\gamma \mathrm{cos}^2\varphi +\mathrm{sin}^2\varphi \right)}}.`$ (12)
The surface integrals (11) and (12) may be evaluated as
$`\mathrm{\Gamma }_{}(\gamma )`$ $`=`$ $`\{\begin{array}{ccc}{\displaystyle \frac{\mathrm{sinh}^1\sqrt{\frac{1\gamma }{\gamma }}}{\left(1\gamma \right)^{\frac{3}{2}}}}{\displaystyle \frac{1}{1\gamma }}\hfill & & \hfill \text{for}\mathrm{\hspace{0.33em}\hspace{0.33em}0}<\gamma <1\\ & & \\ {\displaystyle \frac{1}{\gamma 1}}{\displaystyle \frac{\mathrm{sec}^1\sqrt{\gamma }}{\left(\gamma 1\right)^{\frac{3}{2}}}}\hfill & & \hfill \text{for}\gamma >1\end{array},`$ (16)
$`\mathrm{\Gamma }_{}(\gamma )`$ $`=`$ $`\{\begin{array}{ccc}{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{1}{1\gamma }}{\displaystyle \frac{\gamma \mathrm{sinh}^1\sqrt{\frac{1\gamma }{\gamma }}}{\left(1\gamma \right)^{\frac{3}{2}}}}\right)\hfill & & \hfill \text{for}\mathrm{\hspace{0.33em}\hspace{0.33em}0}<\gamma <1\\ & & \\ {\displaystyle \frac{1}{2}}\left({\displaystyle \frac{\gamma \mathrm{sec}^1\sqrt{\gamma }}{\left(\gamma 1\right)^{\frac{3}{2}}}}{\displaystyle \frac{1}{\gamma 1}}\right)\hfill & & \hfill \text{for}\gamma >1\end{array}.`$ (20)
We exclude the cases of
* the isotropic HCM with $`\gamma =1`$, and
* the anomalous hyperbolic HCM with $`\gamma <0`$
from consideration.
The dyadic Bruggeman equation (4) provides the two nonlinear scalar equations
$`{\displaystyle \frac{ϵ^aϵ_{}^{Br}}{1+D_{}\left(ϵ^aϵ_{}^{Br}\right)}}f_a+{\displaystyle \frac{ϵ^bϵ_{}^{Br}}{1+D_{}\left(ϵ^bϵ_{}^{Br}\right)}}f_b=0,`$ (21)
$`{\displaystyle \frac{ϵ^aϵ_{}^{Br}}{1+D_{}\left(ϵ^aϵ_{}^{Br}\right)}}f_a+{\displaystyle \frac{ϵ^bϵ_{}^{Br}}{1+D_{}\left(ϵ^bϵ_{}^{Br}\right)}}f_b=0,`$ (22)
coupled via $`D_,`$, which can be solved straightforwardly for $`ϵ_{}^{Br}`$ and $`ϵ_{}^{Br}`$ using standard numerical techniques.
## 3 Group velocity
Let us consider a wavepacket which is a superposition of planewaves with phasors
$$\begin{array}{c}\underset{¯}{E}(\underset{¯}{r})=\underset{¯}{E}_0\mathrm{exp}\left(i\underset{¯}{k}\text{ }\text{}\text{ }\underset{¯}{r}\right)\hfill \\ \underset{¯}{H}(\underset{¯}{r})=\underset{¯}{H}_0\mathrm{exp}\left(i\underset{¯}{k}\text{ }\text{}\text{ }\underset{¯}{r}\right)\hfill \end{array}\}.$$
(23)
The group velocity $`\underset{¯}{v}_g`$ of the wavepacket is conventionally defined in terms of the gradient of the angular frequency $`\omega `$ with respect to $`\underset{¯}{k}`$ ; i.e.,
$$\underset{¯}{v}_g=_{\underset{¯}{k}}\omega |_{\omega =\omega (k_{avg})},$$
(24)
where $`k_{avg}`$ denotes the average wavenumber of the wavepacket. Herein we adopt the compact notation
$$_{\underset{¯}{k}}(\frac{}{k_x},\frac{}{k_y},\frac{}{k_z})$$
(25)
for the gradient operator with respect to $`\underset{¯}{k}`$, where $`(k_x,k_y,k_z)`$ is the representation of $`\underset{¯}{k}`$ in terms of its Cartesian components.
In order to calculate the group velocity in the uniaxial dielectric HCM (3), denoted as $`\underset{¯}{v}_g^{Br}`$, we exploit the corresponding planewave dispersion relation as follows. The combination of (3) with the source–free Maxwell curl postulates
$$\begin{array}{c}\times \underset{¯}{E}(\underset{¯}{r})=i\omega \underset{¯}{B}(\underset{¯}{r})\hfill \\ \times \underset{¯}{H}(\underset{¯}{r})=i\omega \underset{¯}{D}(\underset{¯}{r})\hfill \end{array}\},$$
(26)
delivers the vector Helmholtz equation
$$\left[\left(\times \underset{¯}{\underset{¯}{I}}\right)\text{ }\text{}\text{ }\left(\times \underset{¯}{\underset{¯}{I}}\right)\mu _0\omega ^2\underset{¯}{\underset{¯}{ϵ}}^{Br}\right]\text{ }\text{}\text{ }\underset{¯}{E}_0=\underset{¯}{0},$$
(27)
with $`\mu _0`$ being the permeability of free space. The requirement that (27) provide nonzero solutions for the planewave phasors (23) yields the dispersion relation
$$W(\underset{¯}{k},\omega )=0,$$
(28)
wherein the scalar function $`W`$ is defined as
$$W(\underset{¯}{k},\omega )=\left(\underset{¯}{k}\text{ }\text{}\text{ }\underset{¯}{k}ϵ_{}^{Br}\mu _0\omega ^2\right)\left(\underset{¯}{k}\text{ }\text{}\text{ }\underset{¯}{\underset{¯}{ϵ}}^{Br}\text{ }\text{}\text{ }\underset{¯}{k}ϵ_{}^{Br}ϵ_{}^{Br}\mu _0\omega ^2\right).$$
(29)
The dispersion relation (28) admits two wavevector solutions: the ordinary wavevector $`\underset{¯}{k}_{or}`$ and the extraordinary wavector $`\underset{¯}{k}_{ex}`$, satisfying
$$\begin{array}{c}\underset{¯}{k}_{or}\text{ }\text{}\text{ }\underset{¯}{k}_{or}ϵ_{}^{Br}\mu _0\omega ^2=0\hfill \\ \underset{¯}{k}_{ex}\text{ }\text{}\text{ }\underset{¯}{\underset{¯}{ϵ}}^{Br}\text{ }\text{}\text{ }\underset{¯}{k}_{ex}ϵ_{}^{Br}ϵ_{}^{Br}\mu _0\omega ^2=0\hfill \end{array}\}.$$
(30)
We note that the magnitude of the ordinary wavevector is direction–independent, and the ordinary and extraordinary wavevectors coincide when $`\underset{¯}{k}_{ex}`$ is directed along $`\underset{¯}{\overset{^}{c}}`$.
By taking the gradient of the dispersion relation (28) with respect to $`\underset{¯}{k}`$, we find
$$_{\underset{¯}{k}}W+\frac{W}{\omega }_{\underset{¯}{k}}\omega =\underset{¯}{0}.$$
(31)
Hence, the HCM group velocity (24) may be expressed as
$$\underset{¯}{v}_g^{Br}=\frac{1}{W/\omega }_{\underset{¯}{k}}W|_{\omega =\omega (k_{avg})}.$$
(32)
The partial derivative terms involving $`W`$ are found to be
$`_{\underset{¯}{k}}W`$ $`=`$ $`2\left[\left(\underset{¯}{k}\text{ }\text{}\text{ }\underset{¯}{\underset{¯}{ϵ}}^{Br}\text{ }\text{}\text{ }\underset{¯}{k}\omega ^2\mu _0ϵ_{}^{Br}ϵ_{}^{Br}\right)\underset{¯}{k}+\left(\underset{¯}{k}\text{ }\text{}\text{ }\underset{¯}{k}\omega ^2\mu _0ϵ_{}^{Br}\right)\underset{¯}{\underset{¯}{ϵ}}^{Br}\text{ }\text{}\text{ }\underset{¯}{k}\right],`$ (33)
$`{\displaystyle \frac{W}{\omega }}`$ $`=`$ $`\left(\underset{¯}{k}\text{ }\text{}\text{ }\underset{¯}{k}\omega ^2\mu _0ϵ_{}^{Br}\right)\left\{\underset{¯}{k}\text{ }\text{}\text{ }{\displaystyle \frac{d\underset{¯}{\underset{¯}{ϵ}}^{Br}}{d\omega }}\text{ }\text{}\text{ }\underset{¯}{k}\mu _0\omega \left[2ϵ_{}^{Br}ϵ_{}^{Br}+\omega \left({\displaystyle \frac{dϵ_{}^{Br}}{d\omega }}ϵ_{}^{Br}+ϵ_{}^{Br}{\displaystyle \frac{dϵ_{}^{Br}}{d\omega }}\right)\right]\right\}`$ (34)
$`\mu _0\omega \left(2ϵ_{}^{Br}+\omega {\displaystyle \frac{dϵ_{}^{Br}}{d\omega }}\right)\left(\underset{¯}{k}\text{ }\text{}\text{ }\underset{¯}{\underset{¯}{ϵ}}^{Br}\text{ }\text{}\text{ }\underset{¯}{k}ϵ_{}^{Br}ϵ_{}^{Br}\mu _0\omega ^2\right),`$
with
$$\frac{d\underset{¯}{\underset{¯}{ϵ}}^{Br}}{d\omega }=\frac{dϵ_{}^{Br}}{d\omega }\underset{¯}{\underset{¯}{I}}+\left(\frac{dϵ_{}^{Br}}{d\omega }\frac{dϵ_{}^{Br}}{d\omega }\right)\underset{¯}{\overset{^}{c}}\underset{¯}{\overset{^}{c}}.$$
(35)
By virtue of (30), we see that the ordinary and the extraordinary group velocities are given by
$$\underset{¯}{v}_g^{Br}|_{\underset{¯}{k}=\underset{¯}{k}_{or}}=\frac{2}{\omega \mu _0\left(2ϵ_{}^{Br}+\omega {\displaystyle \frac{dϵ_{}^{Br}}{d\omega }}\right)}\underset{¯}{k}_{or}$$
(36)
and
$$\underset{¯}{v}_g^{Br}|_{\underset{¯}{k}=\underset{¯}{k}_{ex}}=\frac{2}{\omega \mu _0\left[2ϵ_{}^{Br}ϵ_{}^{Br}+\omega \left({\displaystyle \frac{dϵ_{}^{Br}}{d\omega }}ϵ_{}^{Br}+ϵ_{}^{Br}{\displaystyle \frac{dϵ_{}^{Br}}{d\omega }}\right)\right]\underset{¯}{k}_{ex}\text{ }\text{}\text{ }{\displaystyle \frac{d\underset{¯}{\underset{¯}{ϵ}}^{Br}}{d\omega }}\text{ }\text{}\text{ }\underset{¯}{k}_{ex}}\underset{¯}{\underset{¯}{ϵ}}^{Br}\text{ }\text{}\text{ }\underset{¯}{k}_{ex},$$
(37)
respectively.
In order to find the derivatives of $`ϵ_{}^{Br}`$ and $`ϵ_{}^{Br}`$ needed to evaluate the group velocities (36) and (37), we have to exploit the Bruggeman equations (21) and (22). As a precursor, let us first note the derivatives of the depolarization dyadic components
$`{\displaystyle \frac{dD_{}}{d\omega }}`$ $`=`$ $`\alpha _{11}{\displaystyle \frac{dϵ_{}^{Br}}{d\omega }}+\alpha _{12}{\displaystyle \frac{dϵ_{}^{Br}}{d\omega }},`$ (38)
$`{\displaystyle \frac{dD_{}}{d\omega }}`$ $`=`$ $`\alpha _{21}{\displaystyle \frac{dϵ_{}^{Br}}{d\omega }}+\alpha _{22}{\displaystyle \frac{dϵ_{}^{Br}}{d\omega }},`$ (39)
with
$`\alpha _{11}`$ $`=`$ $`{\displaystyle \frac{U_{}^2}{U_{}^2ϵ_{}^{Br}ϵ_{}^{Br}}}\left(\mathrm{\Gamma }_{}+\gamma {\displaystyle \frac{d\mathrm{\Gamma }_{}}{d\gamma }}\right){\displaystyle \frac{\gamma \mathrm{\Gamma }_{}}{\left(ϵ_{}^{Br}\right)^2}},`$ (40)
$`\alpha _{12}`$ $`=`$ $`{\displaystyle \frac{U_{}^2}{U_{}^2\left(ϵ_{}^{Br}\right)^2}}\left(\mathrm{\Gamma }_{}+\gamma {\displaystyle \frac{d\mathrm{\Gamma }_{}}{d\gamma }}\right),`$ (41)
$`\alpha _{21}`$ $`=`$ $`\left({\displaystyle \frac{U_{}^2}{U_{}^2\left(ϵ_{}^{Br}\right)^2}}\right){\displaystyle \frac{d\mathrm{\Gamma }_{}}{d\gamma }},`$ (42)
$`\alpha _{22}`$ $`=`$ $`\left({\displaystyle \frac{U_{}^2ϵ_{}^{Br}}{U_{}^2\left(ϵ_{}^{Br}\right)^3}}\right){\displaystyle \frac{d\mathrm{\Gamma }_{}}{d\gamma }}{\displaystyle \frac{\mathrm{\Gamma }_{}}{\left(ϵ_{}^{Br}\right)^2}},`$ (43)
and
$`{\displaystyle \frac{d\mathrm{\Gamma }_{}}{d\gamma }}`$ $`=`$ $`\{\begin{array}{ccc}{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{3\mathrm{sinh}^1\sqrt{\frac{1\gamma }{\gamma }}}{\left(1\gamma \right)^{\frac{5}{2}}}}{\displaystyle \frac{1+2\gamma }{\left(1\gamma \right)^2\gamma }}\right)\hfill & & \hfill \text{for}\mathrm{\hspace{0.33em}\hspace{0.33em}0}<\gamma <1\\ & & \\ {\displaystyle \frac{1}{2}}\left({\displaystyle \frac{1+2\gamma }{\left(\gamma 1\right)^2\gamma }}+{\displaystyle \frac{3\mathrm{sec}^1\sqrt{\gamma }}{\left(\gamma 1\right)^{\frac{5}{2}}}}\right)\hfill & & \hfill \text{for}\gamma >1\end{array},`$ (47)
$`{\displaystyle \frac{d\mathrm{\Gamma }_{}}{d\gamma }}`$ $`=`$ $`\{\begin{array}{ccc}{\displaystyle \frac{1}{4}}\left({\displaystyle \frac{3}{\left(1\gamma \right)^2}}{\displaystyle \frac{\left(2+\gamma \right)\mathrm{sinh}^1\sqrt{\frac{1\gamma }{\gamma }}}{\left(1\gamma \right)^{\frac{5}{2}}}}\right)\hfill & & \hfill \text{for}\mathrm{\hspace{0.33em}\hspace{0.33em}0}<\gamma <1\\ & & \\ {\displaystyle \frac{1}{4}}\left({\displaystyle \frac{\left(2+\gamma \right)\mathrm{sec}^1\sqrt{\gamma }}{\left(\gamma 1\right)^{\frac{5}{2}}}}+{\displaystyle \frac{3}{\left(\gamma 1\right)^2}}\right)\hfill & & \hfill \text{for}\gamma >1\end{array}.`$ (51)
Now we turn to the Bruggeman equations (21) and (22). Their derivatives with respect to $`\omega `$ may be written as
$`\beta _{11}{\displaystyle \frac{dϵ_{}^{Br}}{d\omega }}+\beta _{12}{\displaystyle \frac{dϵ_{}^{Br}}{d\omega }}+\beta _{13}=0,`$ (52)
$`\beta _{21}{\displaystyle \frac{dϵ_{}^{Br}}{d\omega }}+\beta _{22}{\displaystyle \frac{dϵ_{}^{Br}}{d\omega }}+\beta _{23}=0,`$ (53)
with
$`\beta _{11}`$ $`=`$ $`\alpha _{11}\left(ϵ^aϵ_{}^{Br}\right)\left(ϵ^bϵ_{}^{Br}\right)+D_{}\left(2ϵ_{}^{Br}ϵ^aϵ^b\right)1,`$ (54)
$`\beta _{12}`$ $`=`$ $`\alpha _{12}\left(ϵ^aϵ_{}^{Br}\right)\left(ϵ^bϵ_{}^{Br}\right),`$ (55)
$`\beta _{13}`$ $`=`$ $`\left[f_a+D_{}\left(ϵ^bϵ_{}^{Br}\right)\right]{\displaystyle \frac{dϵ^a}{d\omega }}+\left[f_b+D_{}\left(ϵ^aϵ_{}^{Br}\right)\right]{\displaystyle \frac{dϵ^b}{d\omega }},`$ (56)
$`\beta _{21}`$ $`=`$ $`\alpha _{21}\left(ϵ^aϵ_{}^{Br}\right)\left(ϵ^bϵ_{}^{Br}\right),`$ (57)
$`\beta _{22}`$ $`=`$ $`\alpha _{22}\left(ϵ^aϵ_{}^{Br}\right)\left(ϵ^bϵ_{}^{Br}\right)+D_{}\left(2ϵ_{}^{Br}ϵ^aϵ^b\right)1,`$ (58)
$`\beta _{23}`$ $`=`$ $`\left[f_a+D_{}\left(ϵ^bϵ_{}^{Br}\right)\right]{\displaystyle \frac{dϵ^a}{d\omega }}+\left[f_b+D_{}\left(ϵ^aϵ_{}^{Br}\right)\right]{\displaystyle \frac{dϵ^b}{d\omega }}.`$ (59)
The derivatives of $`ϵ_{}^{Br}`$ and $`ϵ_{}^{Br}`$ therefore finally emerge as
$`{\displaystyle \frac{dϵ_{}^{Br}}{d\omega }}`$ $`=`$ $`{\displaystyle \frac{\beta _{12}\beta _{23}\beta _{22}\beta _{13}}{\beta _{11}\beta _{22}\beta _{12}\beta _{21}}},`$ (60)
$`{\displaystyle \frac{dϵ_{}^{Br}}{d\omega }}`$ $`=`$ $`{\displaystyle \frac{\beta _{21}\beta _{13}\beta _{11}\beta _{23}}{\beta _{11}\beta _{22}\beta _{12}\beta _{21}}}.`$ (61)
To summarize, given a uniaxial dielectric HCM with permittivity dyadic $`\underset{¯}{\underset{¯}{ϵ}}^{Br}`$ estimated using the Bruggeman homogenization formalism, the group velocity (24) may be computed using the expression (32), with (33) and (34), wherein the derivatives of $`ϵ_{}^{Br}`$ and $`ϵ_{}^{Br}`$ are provided by (60) and (61).
## 4 Numerical studies
Without loss of generality, let us choose the axis of rotational symmetry of the component spheroids to lie along the $`x`$ axis, i.e., $`\underset{¯}{\overset{^}{c}}=\underset{¯}{\overset{^}{x}}`$. We consider wavevectors lying in the $`xy`$ plane, oriented at an angle $`\theta `$ to the $`x`$ axis. That is, we take
$$\underset{¯}{k}=k\{\mathrm{cos}\theta ,\mathrm{sin}\theta ,0\}.$$
(62)
Thus, the magnitudes $`k=k_{or}|\underset{¯}{k}_{or}|`$ and $`k=k_{ex}|\underset{¯}{k}_{ex}|`$ of the ordinary and extraordinary wavevectors arise from (30) as
$`k_{or}`$ $`=`$ $`\omega \sqrt{\mu _0ϵ_{}^{Br}},`$ (63)
$`k_{ex}`$ $`=`$ $`\omega \sqrt{{\displaystyle \frac{\mu _0ϵ_{}^{Br}ϵ_{}^{Br}}{ϵ_{}^{Br}\mathrm{cos}^2\theta +ϵ_{}^{Br}\mathrm{sin}^2\theta }}}.`$ (64)
Let us explore numerically the enhancement in group velocity that can arise through homogenization, paying special attention to directional effects induced by the shape of the component spheroidal particles. In particular, we choose the component material phase $`a`$ to have a relatively high permittivity $`ϵ^a`$ and a relatively small frequency–dispersion term $`dϵ^a/d\omega `$, compared with the component material phase $`b`$. As representative constitutive parameter values, we set: $`ϵ^a=30ϵ_0`$, $`\left(dϵ^a/d\omega \right)|_{\omega =\omega _o}=6ϵ_0/\omega _o`$, $`ϵ^b=1.2ϵ_0`$ and $`\left(dϵ^b/d\omega \right)|_{\omega =\omega _o}=12ϵ_0/\omega _o`$, where $`ϵ_0`$ is the permittivity of free space.
In Figure 1, the Bruggeman estimates of the HCM permittivity parameters $`ϵ_{}^{Br}`$ and $`ϵ_{}^{Br}`$ are plotted as functions of volume fraction $`f_a`$, for the range of values of $`\rho =U_{}/U_{}`$ shown in Table 1. Clearly, $`ϵ_,^{Br}ϵ^b`$ as $`f_a0`$ and $`ϵ_,^{Br}ϵ^a`$ as $`f_a1`$. We see that $`ϵ_{}^{Br}`$ becomes an increasingly nonlinear function of $`f_a`$ as $`\rho `$ decreases, whereas $`ϵ_{}^{Br}`$ becomes an increasingly nonlinear function of $`f_a`$ as $`\rho `$ increases.
In Figure 2, the magnitude of the group velocity $`v_g^{Br}=|\underset{¯}{v}_g^{Br}|`$ of a wavepacket in the chosen HCM is plotted against volume fraction. The group velocities are calculated with $`\underset{¯}{k}=\underset{¯}{k}_{ex}`$ for $`\theta =0^{},30^{},60^{}`$ and $`90^{}`$. The corresponding graphs for $`180^{}\theta `$ are the same as those for $`\theta `$. Since the ordinary wavevector $`\underset{¯}{k}_{or}=\underset{¯}{k}_{ex}`$ at $`\theta =0^{}`$, the ordinary group velocities for any $`\theta `$ are identical to those provided in Figure 2(a) wherein the results for $`\theta =0^{}`$ are presented. The group velocity magnitudes for the component material phases $`a`$ and $`b`$ are $`v_g^a=0.166c`$ and $`v_g^b=0.152c`$, respectively (as is confirmed in Figure 2 by the group velocity values at $`f_a=1`$ and $`f_a=0`$, respectively), where $`c=1/\sqrt{ϵ_0\mu _0}`$. Hence, for this particular homogenization example, group–velocity enhancement arises when $`v_g^{Br}>\text{max}\{v_g^a,v_g^b\}=0.166c`$. The group–velocity–enhancement region is identified by shading in Figure 2.
It may be discerned from Figure 2(a) that group–velocity enhancement occurs over an increasingly large range of $`f_a`$ values as $`\rho `$ decreases. Furthermore, the degree of enhancement at $`\rho =20`$ is much smaller than it is at $`\rho =0.05`$.
As $`\theta `$ increases, the range of $`f_a`$ values at which group–velocity enhancement occurs progressively decreases for small values of $`\rho `$. In fact, at $`\theta =60^{}`$ there is no longer any enhancement in group velocity for $`\rho =0.05`$. At $`\theta =90^{}`$, the group–velocity enhancement characteristics at low and high values of $`\rho `$ are approximately the reverse of their respective characteristics at $`\theta =0^{}`$. That is, group–velocity enhancement occurs over a wide range of $`f_a`$ values for high values of $`\rho `$ at $`\theta =90^{}`$, but there is no enhancement in group velocity at low values of $`\rho `$.
Clearly therefore, enhancement of group velocity is maximum in a direction parallel to the longest semi–axis of the spheroidal particles, which can be prolate ($`\rho <1`$) or oblate ($`\rho >1)`$. For spherical particles ($`\rho =1`$), group–velocity enhancement is direction–independent, and we recover the results of the predecessor study .
## 5 Concluding remarks
The enhancement in group velocity brought about by homogenization is sensitively dependent upon directional properties. Both the shape of the component spheroidal particles, and their orientation relative to the direction of propagation, strongly influence the group–velocity enhancement.
The homogenization scenario presented here deals with the conceptualization of a uniaxial HCM as arising from identically oriented spheroidal particles of isotropic component material phases. The homogenization of two uniaxial dielectric component phases distributed as spherical particles is mathematically equivalent, provided that the distinguished axes of the component material phases have the same orientation .
Acknowledgement: We thank Prof. G.W. Milton for drawing our attention to group–velocity enhancement in isotropic dielectric composite mediums with HCM permittivity estimated using the Maxwell Garnett formalism, as described in Refs. and .
Table 1. Key for the values of $`\rho =U_{}/U_{}`$ used in Figures 1 and 2. |
warning/0506/hep-ph0506115.html | ar5iv | text | # 𝐵_𝑠→𝜇⁺𝜇⁻ and the upward-going muon flux from the WIMP annihilation in the sun or the earth
## 1 Introduction
There is now compelling evidence for a non-baryonic cold dark matter (DM) component in the universe . In the Minimal Supersymmetric Standard Model (MSSM) with $`R`$ parity, the lightest supersymmetric particle (LSP) is stable and becomes a good candidate for cold dark matter in the universe . The LSP is often the lightest neutralino which is the admixture of Bino, Wino and Higgsinos in the MSSM. In this case, the neutralino DM in our galactic halo might be detected via its elastic scattering with terrestrial nuclear targets . In fact, the DAMA Collaboration even claimed an evidence for DM. However, the CDMS II experiment has reported the upper limit on the DM scattering cross section, which is not compatible with the results of the DAMA experiment. There are several experiments going on searching for the DM scattering at the level of $`\sigma _{\chi p}^{\mathrm{SI}}10^7`$ pb or less. In the most widely studied minimal supergravity (mSUGRA) scenario (or the constrained MSSM), the spin-independent neutralino-proton scattering cross section $`\sigma _{\chi p}`$ turns out very small ($`10^8`$ pb). However there is no solid theoretical rationale for the minimal supergravity scenario, and it is important to calculate the possible maximal values for $`\sigma _{\chi p}`$ in general supergravity scenarios beyond the mSUGRA scenario. And it is very important to impose all the relevant constraints from various experiments in order not to overestimate the cross section. Some important constraints include the lower bounds on the Higgs and SUSY particle masses, $`BX_s\gamma `$ branching ratio, the muon $`(g2)_\mu `$, etc.. One may also take some theoretical consideration on the absence of the color-chrage breaking minima or the directions unbounded from below, etc..
In a previous work , we pointed out that there is a strong correlation between the spin independent neutralino-proton scattering cross section $`\sigma _{\chi p}`$ and the branching ratio for $`B_s\mu ^+\mu ^{}`$ decay within mSUGRA and its extensions. The origin of this correlation resides in the dependence of both observables on $`\mathrm{tan}\beta `$ and the neutral Higgs boson masses; both observables increase for large $`\mathrm{tan}\beta `$ and low Higgs masses. In particular, we have shown that the current upper limit on $`B(B_s\mu ^+\mu ^{})`$ excludes substantial parameter space where the DM scattering cross section is within the CDMS sensitivity region (see also for a detailed analysis).
In this work, we extend our previous study to the indirect detection of neutralino DM with neutrino telescope through upward-going muon flux, and its correlation with $`B(B_s\mu ^+\mu ^{})`$. The energetic neutrino(-induced muon) flux from neutralino DM annihilation in the sun and the earth is one of the promising signals in the indirect detection of neutralino DM . Neutralino DM particles in the halo can be captured by the sun or by the earth, when their velocities drop below escape velocities via their elastic scattering with matter in the sun or earth. Then they will accumulate in the core of the sun and the earth and will eventually annihilate into ordinary SM particles. Among the annihilaion products, neutrinos can pass through the sun and the earth, and then could be detected in neutrino telescopes through their conversion to muons via charged-current scattering with neuclei near the detectors. Baksan , MACRO , Super-K and AMANDA released upper limits on the upward-going muon flux. There are also planned or proposed neutrino telescopes such as ANTARES , IceCube and NESTOR etc..
An important point of the indirect detection of neutralino DM with neutrino telescopes is that the neutrino flux strongly depends on the capture rate of neutralino by the sun or the earth, which in turn depends on neutralino-nucleon scattering cross sections. Therefore we expect some correlation between the neutrino flux and $`B(B_s\mu ^+\mu ^{})`$, which is similar to the strong correlation between $`\sigma _{\chi p}`$ and $`B(B_s\mu ^+\mu ^{})`$ as discussed in Ref. . Indeed, we will show that the current upper limit of $`B(B_s\mu ^+\mu ^{})<4.1\times 10^7`$ (90 % CL) puts strong constraints on the upward-going muon flux in the supersymmetric models which give rather large spin-independent neutralino-proton scattering cross section.
This paper is organized as following. In Sec. 2, we give a brief review on the indirect detection of the DM through the upward-going muon flux. In Sec. 3, we consider the upward-going muon fluxes in some supergravity scenarios and illuminate our point that $`B_s\mu ^+\mu ^{}`$ branching ratio plays an important role. In Sec. 4, we summarize the results.
## 2 Indirect detection through the upward-going muon flux
As we mentioned in the introdecution, the observation of energetic neutrinos from the sun and/or the earth would provide convincing evidence of the existence of neutralino dark matter in galactic halo . The flux of energetic neutrinos from neutralino annihilation in the sun or the earth is proportional to the rate of neutralino annihilation in the sun or in the earth and the energy spectrum of neutrinos from the annihilation. The time evolution of the number of neutralino, $`N`$ in the sun (or in the earth) is given by
$`\dot{N}=CC_AN^2`$ (1)
where $`C`$ is the capture rate of neutralino by the sun or the earth and $`C_A`$ is the total annihilation cross section times relative velocity per volume. From Eq.(2.1), we find that the present annihilation rate is
$`\mathrm{\Gamma }_A={\displaystyle \frac{1}{2}}C_AN^2={\displaystyle \frac{1}{2}}C\mathrm{tanh}^2(\sqrt{CC_A}t_0)`$ (2)
where $`t_04.5`$ Gyr is the age of the solar system. For $`\sqrt{CC_A}t_01`$, the annihilation rate is $`\mathrm{\Gamma }_A\frac{1}{2}C^2C_At_0^2`$ and less than its maximal value. But, for $`\sqrt{CC_A}t_01`$, the neutralino density reach equilibrium and the annihilation rate is $`\mathrm{\Gamma }_A\frac{1}{2}C`$. Therefore, when accretion is efficient, the annihilation rate depends on the capture rate $`C`$, but not on the annihilation cross section.
In turn, the capture rate $`C`$ strongly depends on the elastic scattering cross section of neutralino with matter in the sun and the earth. The capture rate for the earth primarily depends on the spin-independent DM scattering cross section. For the capture rate in the sun, however, both spin-independent and spin-dependent scattering cross section can be important and the significance of each contribution depends on the specific SUSY scenarios.
In MSSM, $`t`$-channel Higgs boson and $`s`$-channel squark exchange processes contribute to the spin-independent (scalar) scattering between neutralino and quarks. In many case, dominant contribution to the scalar cross section comes from the Higgs exchange process, which increases for large $`\mathrm{tan}\beta `$ and small Higgs masses and also if neutralino is a mixed gaugino-Higgsino state. On the other hand, for the spin-dependent cross section, $`t`$-chennel $`Z`$ boson and $`s`$-chennel squark exchange processes contribute. Usually $`Z`$ exchange contribution dominates, which is sensitive to Higgsino components of LSP, but largely independent of $`\mathrm{tan}\beta `$. Note that if the Higgsino component of the LSP increases, then both the spin-independent and the spin-dependent scattering cross sections will be enhanced, as shown below in the nonuniversal Higgs mass parameter case.
The capture rate $`C`$ also depends on the local density of neutralino, $`\rho _\chi `$ and the neutralino velocity dispersion in the halo, $`\overline{v}`$ etc. For our numerical calculation, we use the code DARKSUSY and fix $`\overline{v}=270\mathrm{km}/\mathrm{s}`$. For the local density of neutralino we fix $`\rho _\chi =0.3\mathrm{GeV}/\mathrm{cm}^3`$ if $`\mathrm{\Omega }_\chi h^20.025`$, while performing a rescaling of the density as $`\rho _\chi \rho _\chi (\mathrm{\Omega }_\chi h^2/0.025)`$ if $`\mathrm{\Omega }_\chi h^2<0.025`$. Here $`\mathrm{\Omega }_\chi `$ is the neutralino relic density in units of the critical density and $`h`$ is the present Hubble constant in units of $`100\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$.
## 3 Upward-going muon flux in SUSY models
### 3.1 mSUGRA
In mSUGRA model, we assume a universal SUSY breaking scalar mass $`m`$, a universal gaugino mass $`M`$ and a universal trilinear coupling $`A`$ at GUT scale $`m_{\mathrm{GUT}}2\times 10^{16}`$ GeV. We also require that electroweak symmetry break radiatively and then the Higgsino mass parameter $`\mu `$ is determined by the condition:
$$\mu ^2=\frac{m_{H_d}^2m_{H_u}^2\mathrm{tan}^2\beta }{\mathrm{tan}^2\beta 1}\frac{1}{2}M_Z^2.$$
(3)
where $`\mathrm{tan}\beta `$ is the ratio of two Higgs vacuum expectation values and $`m_{H_u}^2,m_{H_d}^2`$ are the soft breaking Higgs masses-squared. With the above mSUGRA assumptions, $`|\mu |`$ is usually large so that the lightest neutralino is bino-like and the pseudo-scalar Higgs mass $`m_A`$ is rather large.
In Fig.s 1 (a) and (b), we show the allowed ranges of the upward-going muon fluxes from the sun and the earth respectively as functions of the LSP mass. The three branches correspond to $`\mathrm{tan}\beta =10,35`$ and 50 cases (from the bottom to the top), respectively. Here, we took $`A=0`$ and $`\mu >0`$ (motivated by the muon $`(g2)_\mu `$ experiment) and varied $`m`$ and $`M`$ up to 1 TeV. We have imposed the experimental bounds for the Higgs and sparticle masses and for $`bs\gamma `$ branching ratio. We also required that the lightest neutralino is LSP. For opposite sign of $`\mu `$, the muon flux could be smaller. But we are interested in the possible maximal values, and we consider the positive $`\mu `$ case only in this work.
In our scan, the muon flux from the sun reaches up to $`20km^2yr^1`$ if the neutralino LSP is light enough ($`m_\chi 100`$ GeV). In the small $`m_\chi `$ region, the neutralino density in the sun can reach (near) equilibrium so that the muon flux is more or less determined by the capture rate of neutralino by the sun. And in turn, the capture rate in the sun is determined primarily by the spin-dependent scattering cross section. (though the contribution from the spin-independent scattering cross section to the capture rate can be comparable to the one from the spin-dependent cross section for very large $`\mathrm{tan}\beta `$ cases) As we already mentioned in Sec. 2, the spin-dependent scattering cross section is largely independent of $`\mathrm{tan}\beta `$. Therefore the upward-going muon flux from the sun in the small $`m_\chi `$ region gives similar values for the three choices of $`\mathrm{tan}\beta `$ values, as one can check from the figure.
In large $`m_\chi `$ region, the neutralino density is usually far from equilibrium since the elastic scattering cross section of the DM with ordinary matter in the sun becomes smaller, and the neutralino annhilation cross section becomes important for the prediction of the muon flux. An important process in this case is the neutralino pair annihilation into $`b\overline{b}`$ through s-channel pseudo-scalar Higgs exchange diagram, which is strongly enhanced for large $`\mathrm{tan}\beta `$ . From the Fig. 1 (a), one can notice a clear dependence of the muon flux from the sun on $`\mathrm{tan}\beta `$ in the large $`m_\chi `$ region.
For the muon flux from the earth, the neutralino density is far less than the equilibrium values and both the capture rate and the annhilation rate are important for the calculation of the muon flux. The resulting flux is much below the one from the sun. The maximal value of the muon flux from the earth is about $`3\times 10^5km^2yr^1`$, which is far below the SUPER-K and AMANDA II sensitivity regions.
Note that there is no further constraint from the $`B_s\mu ^+\mu ^{}`$ bound for the mSUGRA case, once we impose the constraints from the lower bounds for Higgs boson and SUSY particle masses and the $`BX_s\gamma `$ branching ratio, as discussed in Ref. .
### 3.2 Non-universal Higgs model (NUHM)
In the previous subsection, we have shown that the mSUGRA assumption predicts the muon fluxes from the sun and the earth that are far below the sensitivity region of the current experiments. This is mainly because the lightest neutralino is bino-like and the pseudo-scalar Higgs mass is large in the scanned region of mSUGRA scenario. Larger muon flux from the sun and the earth can be obtained if we relax the universal boundary condition at GUT scale.
In this subsection we consider the non-universal Higgs model, in which the assumption of universal soft scalar masses are relaxed for soft Higgs masses, as follows:
$$m_{H_u}^2=m^2(1+\delta _{H_u}),m_{H_d}^2=m^2(1+\delta _{H_d}),$$
(4)
whereas other scalar masses still have a universal mass $`m`$ at GUT scale. Here $`\delta `$’s are parameters with $`O(1)`$.
As an optimal choice for enhancing the muon flux from the sun and the earth, we take the numerical values of $`\delta ^{}`$s as $`\delta _{H_d}=1`$ and $`\delta _{H_u}=1`$. For the postive $`\delta _{H_u}`$, $`\mu `$ becomes lower and the Higgsino component in the neutralino LSP increases so that $`\sigma _{\chi p}`$ is enhanced, as discussed in Ref. . The change of $`|\mu |`$ also has an impact on the Higgs masses because
$$m_A^2=m_{H_u}^2+m_{H_d}^2+2\mu ^2m_{H_d}^2+\mu ^2M_Z^2/2$$
at weak scale. For the negative $`\delta _{H_d}`$, $`m_A`$ and $`m_H`$ become further lower. As the result, both spin-independent scattering cross sections $`\sigma _{\chi p}^{scalar}`$ and spin-dependent one $`\sigma _{\chi p}^{spin}`$ are enhanced compared to mSUGRA case.
In Fig. 2, we present $`\sigma _{\chi p}^{scalar}`$ vs. $`\sigma _{\chi p}^{spin}/(2m_\chi /GeV)`$ in the NUHM (black points) and mSUGRA scenario (green points) for (a) $`\mathrm{tan}\beta =35`$ and (b) $`\mathrm{tan}\beta =50`$ respectively. The dashed straight line in the figure indicates the region in which the two contributions to the capture rate are similar to each other. We observe that both $`\sigma _{\chi p}^{spin}`$ and especially $`\sigma _{\chi p}^{scalar}`$ in NUHM are enhanced a lot compared to mSUGRA scenario. An important point we notice from the figure is that $`\sigma _{\chi p}^{scalar}`$ is usually (especially in the region of large cross section) larger than $`\sigma _{\chi p}^{spin}/(2m_\chi /GeV)`$ in the NUHM case, while the opposite is true for the mSUGRA case. This fact implies the capture rate for the sun in the NUHM is largely determined by the spin-independent scattering cross section rather than spin-dependent one, unlike the mSUGRA scenario. This is because the ratio of the contribution from spin-independent and spin-dependent cross section to the capture rate for the sun is approximately proportional to the ratio of $`\sigma _{\chi p}^{scalar}`$ and $`\sigma _{\chi p}^{spin}/(2m_\chi /GeV)`$.
As we have shown in the previous paper , the current experimental limit of $`B(B_s\mu ^+\mu ^{})`$ puts a strong constraint on the allowed range of the spin-independent cross section. Since the muon flux from the sun and the earth strongly depends on the spin-independent cross section, we naturally expect that the current limit of $`B(B_s\mu ^+\mu ^{})`$ play an important part in restricting the muon flux. This point can be observed clearly in Fig. 3, where we show explicitly the correlation between $`B(B_s\mu ^+\mu ^{})`$ and the muon flux from the sun (a) and the earth (b) in NUHM for $`\mathrm{tan}\beta =35`$ and 50. Note that the $`B_s\mu ^+\mu ^{}`$ is stronger for larger $`\mathrm{tan}\beta `$, and the resulting muon flux becomes smaller for the larger $`\mathrm{tan}\beta `$ case, like the spin independent DM scattering cross section .
The enhancements of the neutralino DM scattering cross sections (both spin-dependent and spin-independent) lead to the substantial change of the muon flux both from the sun and the earth compared with the mSUGRA case. Fig.s 4 (a) and (b) show the muon fluxes from the sun and the earth, respectively, in non-universal Higgs mass scenario with $`\delta _{H_d}=1,\delta _{H_u}=+1`$ for $`\mathrm{tan}\beta =35`$ case. Now the maximal values of the muon fluxes from the sun and the earth are $`10^3(10)km^2yr^1`$, which is two (eight) orders of magnitude lager than the one for the mSUGRA case with $`\mathrm{tan}\beta =35`$. In Fig. 4 (c) and (d), we show the muon flux from the sun and the earth, respectively, in non-universal Higgs mass scenario with $`\delta _{H_d}=1,\delta _{H_u}=+1`$ for $`\mathrm{tan}\beta =50`$. The red points (the open circles) are exclued by the current upper limit of $`B(B_s\mu ^+\mu ^{})`$.
### 3.3 D-brane model
Next, we consider a specific $`D`$ brane model where the SM gauge groups and 3 generations live on different $`Dp`$ branes . In this model, scalar fermion masses are not completely universal and gaugino mass unification can be relaxed. Also the string scale is around $`10^{12}`$ GeV (the intermediate scale) rather than GUT scale.
Since there are now three moduli ($`T_i`$) and one dilaton superfields in this case, we use the following parametrization that is appropriate for several $`T_i`$ moduli:
$`F^S`$ $`=`$ $`\sqrt{3}(S+S^{})m_{3/2}\mathrm{sin}\theta ,`$
$`F^i`$ $`=`$ $`\sqrt{3}(T_i+T_i^{})m_{3/2}\mathrm{cos}\theta \mathrm{\Theta }_i`$ (5)
where $`\theta `$ and $`\mathrm{\Theta }_i(i=1,2,3)`$ with $`_i|\mathrm{\Theta }_i|^2=1`$ parametrize the directions of the goldstinos in the $`S,T_i`$ field space. Then, the gaugino masses are given by
$`M_3`$ $`=`$ $`\sqrt{3}m_{3/2}\mathrm{sin}\theta ,`$
$`M_2`$ $`=`$ $`\sqrt{3}m_{3/2}\mathrm{\Theta }_1\mathrm{cos}\theta ,`$
$`M_Y`$ $`=`$ $`\sqrt{3}m_{3/2}\alpha _Y(M_I)\left({\displaystyle \frac{2\mathrm{\Theta }_3\mathrm{cos}\theta }{\alpha _1(M_I)}}+{\displaystyle \frac{\mathrm{\Theta }_1\mathrm{cos}\theta }{\alpha _2(M_I)}}+{\displaystyle \frac{2\mathrm{sin}\theta }{3\alpha _3(M_I)}}\right),`$ (6)
where
$$\frac{1}{\alpha _Y(M_I)}=\frac{2}{\alpha _1(M_I)}+\frac{1}{\alpha _2(M_I)}+\frac{2}{3\alpha _3(M_I)}.$$
(7)
The string scale $`M_I`$ is determined to be $`M_I=10^{12}(5\times 10^{14})`$ GeV from the $`U(1)_1`$ gauge coupling $`\alpha _1(M_I)=0.1(1)`$ . Note that the gaugino masses are non universal in a natural way in this scenario, unlike other scenarios studied in the previous subsections.
The soft masses for the sfermions and Higgs fields are given by
$`m_Q^2`$ $`=`$ $`m_{3/2}^2\left[1{\displaystyle \frac{3}{2}}\left(1\mathrm{\Theta }_1^2\right)\mathrm{cos}^2\theta \right],`$
$`m_{u^c}^2`$ $`=`$ $`m_{3/2}^2\left[1{\displaystyle \frac{3}{2}}\left(1\mathrm{\Theta }_3^2\right)\mathrm{cos}^2\theta \right],`$
$`m_{d^c}^2`$ $`=`$ $`m_{3/2}^2\left[1{\displaystyle \frac{3}{2}}\left(1\mathrm{\Theta }_2^2\right)\mathrm{cos}^2\theta \right],`$
$`m_L^2`$ $`=`$ $`m_{3/2}^2\left[1{\displaystyle \frac{3}{2}}\left(\mathrm{sin}^2\theta +\mathrm{\Theta }_3^2\mathrm{cos}^2\theta \right)\right],`$
$`m_{e^c}^2`$ $`=`$ $`m_{3/2}^2\left[1{\displaystyle \frac{3}{2}}\left(\mathrm{sin}^2\theta +\mathrm{\Theta }_1^2\mathrm{cos}^2\theta \right)\right],`$
$`m_{H_2}^2`$ $`=`$ $`m_{3/2}^2\left[1{\displaystyle \frac{3}{2}}\left(\mathrm{sin}^2\theta +\mathrm{\Theta }_2^2\mathrm{cos}^2\theta \right)\right],`$
$`m_{H_1}^2`$ $`=`$ $`m_L^2.`$ (8)
Note that the scalar mass universality in the sfermion masses and Higgs masses is achieved when
$$\mathrm{sin}^2\theta =\frac{1}{4}\mathrm{and}\mathrm{\Theta }_i^2=\frac{1}{3}\mathrm{for}i=1,2,3.$$
(9)
And in this case the gaugino masses becomes also universal, when we take only positive numbers for the solutions. For other choices of goldstino angles, the scalar and the gaugino masses become nonuniversal, and there could be larger or smaller flavor violation in the low energy processes as well as enhanced SUSY contributions to the $`a_\mu ^{\mathrm{SUSY}}`$.
The trilinear couplings are given by
$`A_u`$ $`=`$ $`{\displaystyle \frac{\sqrt{3}}{2}}m_{3/2}\left[(\mathrm{\Theta }_2\mathrm{\Theta }_1\mathrm{\Theta }_3)\mathrm{cos}\theta \mathrm{sin}\theta \right],`$
$`A_d`$ $`=`$ $`{\displaystyle \frac{\sqrt{3}}{2}}m_{3/2}\left[(\mathrm{\Theta }_3\mathrm{\Theta }_1\mathrm{\Theta }_2)\mathrm{cos}\theta \mathrm{sin}\theta \right],`$
$`A_e`$ $`=`$ $`0.`$ (10)
Therefore the $`D`$ brane model considered in this work is specified by following six parameters :
$$m_{3/2},\mathrm{tan}\beta ,\theta ,\mathrm{\Theta }_{i=1,2},\mathrm{sign}(\mu ).$$
Due to the departure from the universlity of scalar masses and the proportionality of trilinear couplings, the flavor violation could be different from the mSUGRA case. For example, it is possible to have smaller $`bs`$ transition due to the smaller $`\stackrel{~}{t}_L\stackrel{~}{t}_R`$ mixing and larger stop masses in this $`D`$brane scenarios, so that the $`B(B_s\mu ^+\mu ^{})`$ constraint can be relaxed. This can be seen in Fig. 5 (a) and (b), where the large flux signals are excluded by Super-K and AMANDA II, but not by the $`B(B_s\mu ^+\mu ^{})`$ constraint. In this limited parameter space, one can have a large DM scattering cross section and the upward-going muon flux without conflict with the $`B\mu ^+\mu ^{}`$ branching ratio. Also there is no strong correlations among these observables. Therefore the indirect search for the DM annihilation is complementary to the $`B_s\mu ^+\mu ^{}`$ branching ratio in the $`D`$brane scenarios.
## 4 Conclusions
In this work, we considered the indirect detection of the DM through the upward-going muon flux from the DM annihilation at the core of the sun or the earth, along with the upper bound on the branching ratio for the $`B_s\mu ^+\mu ^{}`$ decay, in some general supergravity scenarios where the upward-going muon flux could be enhanced very much compared to the mSUGRA case. In general supergravity scenario with non-universal Higgs model, we found the following:
* Both $`\sigma _{\chi p}^{spin}`$ and $`\sigma _{\chi p}^{scalar}`$ can be enhanced a lot compared to the mSUGRA scenario, but the enhancement in the spin-independent part is much greater.
* Therefore, contrary to the usual claim, the upward-going muon flux from the sun can be dominated by the spin-independent part $`\sigma _{\chi p}^{scalar}`$ in the NUHM, rahter than by the spin-dependent part $`\sigma _{\chi p}^{spin}`$, as in the mSUGRA scenario \[ Fig. 2 (a) and (b) \].
* The current upper bound $`B(B_s\mu ^+\mu ^{})<4.1\times 10^7`$ excludes a large parameter space where the muon fluxes could be enhanced otherwise, and the constraint is stronger for larger $`\mathrm{tan}\beta `$ \[ Fig. 3 (a) and (b) \].
* The upper bound on $`B(B_s\mu ^+\mu ^{})`$ becomes much stronger than the upper limits on the muon flux from Super-K and AMANDA II \[ Fig. 4 (a)–(d) \].
In the $`D`$brane models with nonuniversal scalar fermion masses, the correlations between the muon flux and $`B(B_s\mu ^+\mu ^{})`$ becomes lost, and the upper bound on $`B(B_s\mu ^+\mu ^{})`$ is complementary to the upper bounds on the muon fluxes from Super-K and AMANDA II. Our study shows that the muon flux originated from the DM annihilation in the sun could be in the range of a few $`\times 10^3`$ /km$`{}_{}{}^{2}`$ yr.
Our study indicates that it is most important to include the $`B_s\mu ^+\mu ^{}`$ branching ratio constraint when we study the direct and the indirect detections of the neutralino DM in general supergravity scenarios. The upper limit on the $`B_s\mu ^+\mu ^{}`$ branching ratio excludes significant part of parameter space where the DM scattering cross section and the upward-going muon flux could be enhanced above/around the current experiments. Unless the chargino-stop contribution to $`B_s\mu ^+\mu ^{}`$ is very small or there is fortuitous cancellation between the chargino-stop and the gluino-sbottom loop contributions, the spin-independent DM scattering cross section and the indirect detection rate through the upward-going muon flux are strongly contrained by the $`B_s\mu ^+\mu ^{}`$ branching ratio. Since both the direct and the indirect detection rates are well below the current experiments in most supergravity model parameter space when the $`B_s\mu ^+\mu ^{}`$ branching ratio constraint is imposed, it would be a great challenge for experimemtalists to reach such sensitivity to have positive signals of the DM search.
###### Acknowledgments.
The authors are grateful to A. Masiero and P. Ullio for useful discussions on indirect detection of the neutralino DM. PK is supported in part by KOSEF Sundo grant R02-2003-000-10085-0, KOSEF through CHEP at Kyungpook National University and KRF grant KRF-2002-070-C00022. The work of YGK was supported by Korea Research Foundation and the Korean Federation of Science and Technology Societies Grant funded by Korea Government (MOEHRD, Basic Research Promotion Fund). |
warning/0506/math0506579.html | ar5iv | text | # Semi-direct products of Lie algebras and their invariants
## Introduction
The ground field 𝕜 is algebraically closed and of characteristic zero. The goal of this paper is to extend the standard invariant-theoretic design, well-developed in the reductive case, to the setting of non-reductive group representations. This concerns the following notions and results: the existence of generic stabilisers and generic isotropy groups for (finite-dimensional rational) representations; structure of the fields and algebras of invariants; quotient morphisms and structure of their fibres. One of the main tools for obtaining non-reductive Lie algebras is the semi-direct product construction. There is a number of articles devoted to the study of the coadjoint representations of non-reductive Lie algebras; in particular, semi-direct products, see e.g. . In this article, we consider such algebras from a broader point of view. In particular, we found that the adjoint representation is an interesting object, too. Our main references for Invariant Theory are and . All algebraic groups are assumed to be linear.
If an algebraic group $`A`$ acts on an affine variety $`X`$, then $`\text{𝕜}[X]^A`$ stands for the algebra of $`A`$-invariant regular functions on $`X`$. If $`\text{𝕜}[X]^A`$ is finitely generated, then $`X//A:=\mathrm{Spec}\text{𝕜}[X]^A`$, and the quotient morphism $`\pi _A:XX//A`$ is the mapping associated with the embedding $`\text{𝕜}[X]^A\text{𝕜}[X]`$. If $`\text{𝕜}[X]^A`$ is polynomial, then the elements of any set of algebraically independent homogeneous generators will be referred to as basic invariants.
Let $`G`$ be a connected reductive algebraic group with Lie algebra $`g`$. Choose a Cartan subalgebra $`tg`$ with the corresponding Weyl group $`𝖶`$. The adjoint representation $`(G:g)`$ has a number of good properties, some of which are listed below:
$``$ The adjoint representation is self-dual, and $`t`$ is a generic stabiliser for it;
$``$ The algebra of invariants $`\text{𝕜}[g]^G`$ is polynomial;
$``$ the restriction homomorphism $`\text{𝕜}[g]\text{𝕜}[t]`$ induces the isomorphism $`\text{𝕜}[g]^G\text{𝕜}[t]^𝖶`$ (Chevalley’s theorem);
$``$ The quotient morphism $`\pi _G:gg//G`$ is equidimensional and the fibre of the origin, $`𝒩:=\pi _G^1(\pi _G(0))`$, is an irreducible complete intersection. The ideal of $`𝒩`$ in $`\text{𝕜}[g]`$ is generated by the basic invariants;
$``$ $`𝒩`$ is the union of finitely many $`G`$-orbits.
Each of these properties may fail if $`g`$ is replaced with an arbitrary algebraic Lie algebra $`q`$. In particular, one have to distinguish the adjoint and coadjoint representations of $`q`$. As usual, $`\mathrm{ad}`$ (resp. $`\mathrm{ad}^{}`$) stands for the adjoint (resp. coadjoint) representation. Write $`Q`$ for a connected group with Lie algebra $`q`$.
First, we consider the problem of existence of generic stabilisers for $`\mathrm{ad}`$ and $`\mathrm{ad}^{}`$. (See Section 1 for precise definitions). It turns out that if $`(q,\mathrm{ad})`$ has a generic stabiliser, say $`h`$, then $`h`$ is commutative and $`n_q(h)=h`$. This yields a Chevalley-type theorem for the fields of invariants: $`\text{𝕜}(q)^Q\text{𝕜}(h)^W`$, where $`W=N_Q(h)/Z_Q(h)`$ is finite. We also notice that $`(q,\mathrm{ad})`$ has a generic stabiliser if and only if the Cartan subalgebras of $`q`$ are commutative. If $`(q,\mathrm{ad}^{})`$ has a generic stabiliser, say $`h`$, then $`h`$ is commutative, $`dimN_Q(h)=dim(q^{})^h`$, and $`\text{𝕜}(q^{})^Q\text{𝕜}((q^{})^h)^{N_Q(h)}`$. But unlike the adjoint case, the action $`(N_Q(h):(q^{})^h)`$ does not necessarily reduce to a finite group action. We prove that under a natural constraint the representation of the identity component of $`N_Q(h)`$ on $`(q^{})^h`$ is the coadjoint representation.
Our main efforts are connected with the following situation. Suppose that $`(q,\mathrm{ad})`$ or $`(q,\mathrm{ad}^{})`$ has some of the above good properties and $`V`$ is a (finite-dimensional rational) $`Q`$-module. Form the Lie algebra $`qV`$. It is the semi-direct product of $`q`$ and $`V`$, $`V`$ being a commutative ideal in it. The corresponding connected algebraic group is $`QV`$. (See section 4 for the details.) Then we want to realise to which extent those good properties are preserved under this procedure. This surely depends on $`V`$, and we are essentially interested in two cases:
(a) $`q`$ is arbitrary and $`V=q`$ or $`q^{}`$ (the adjoint or coadjoint $`q`$-module);
(b) $`q=g`$ is reductive and $`V`$ is an arbitrary $`G`$-module.
For (a), we prove that if $`(q,\mathrm{ad})`$ has a generic stabiliser, then so do $`(qq,\mathrm{ad})`$ and $`(qq^{},\mathrm{ad})`$. Furthermore, the passages $`qqq`$ and $`qqq^{}`$ does not affect the generalised Weyl group $`W`$, and both fields $`\text{𝕜}(qq)^{Qq}`$ and $`\text{𝕜}(qq^{})^{Qq^{}}`$ are purely transcendental extensions of $`\text{𝕜}(q)^Q`$. It is also true that if $`(q,\mathrm{ad}^{})`$ has a generic stabiliser, then so does $`(qq,\mathrm{ad}^{})`$.
For (b), we prove that $`(gV,\mathrm{ad})`$ always has a generic stabiliser. But this is not the case for $`\mathrm{ad}^{}`$. Recall that any $`g`$-module $`V`$ has a generic stabiliser. The following result seems to be quite unexpected. Suppose generic $`G`$-orbits in $`V`$ are closed (i.e., the action $`(G:V)`$ is stable), then $`(gV,\mathrm{ad}^{})`$ has a generic stabiliser if and only if $`V`$ is a polar $`G`$-module in the sense of . The assumption of stability is relatively harmless, since there are only finitely many $`G`$-modules without that property. On the other hand, the hypothesis of being polar is quite restrictive, because for any $`G`$ there are only finitely many polar representations.
One of our main observations is that there are surprisingly many nonreductive Lie algebras $`a`$ and $`a`$-modules $`M`$ such that $`\text{𝕜}[M]^A`$ is a polynomial algebra. Furthermore, the basic invariants of $`\text{𝕜}[M]^A`$ can explicitly be constructed using certain modules of covariants. This concerns the following cases:
– If $`g`$ is reductive and $`V`$ is an arbitrary $`g`$-module, then one takes $`a=M=gV`$;
– If the action $`(Q:V)`$ satisfies some good properties, then one takes $`a=qq`$ and $`M=VV`$. Furthermore, the passage $`(q,V)(\widehat{q}=qq,\widehat{V}=VV)`$ can be iterated.
The precise statements are given below.
0.1 Theorem. Let $`V`$ be an arbitrary $`G`$-module. Set $`q=gV`$, $`Q=GV`$, and $`m=dimV^t`$. Notice that $`1V`$ is a commutative normal subgroup of $`Q`$ (in fact, the unipotent radical of $`Q`$). Then
* $`\text{𝕜}[q]^{1V}`$ is a polynomial algebra of Krull dimension $`dimg+m`$. It is freely generated by the coordinates on $`g`$ and the functions $`\widehat{F}_i`$, $`i=1,\mathrm{},m`$, associated with covariants of type $`V^{}`$.
* $`\text{𝕜}[q]^Q`$ is a polynomial algebra of Krull dimension $`dimt+m`$. It is freely generated by the basic invariants of $`\text{𝕜}[g]^G`$ and the same functions $`\widehat{F}_i`$, $`i=1,\mathrm{},m`$.
* $`\mathrm{max}dim_{xq}Qx=dimqdimq//Q`$;
* If $`\pi :qq//Q`$ is the quotient morphism and $`\mathrm{\Omega }:=\{xq\text{d}\pi _x\text{ is onto }\}`$, then $`q\mathrm{\Omega }`$ contains no divisors.
Given a $`q`$-module $`V`$, the space $`V\times V`$ can be regarded as $`qq`$-module in a very natural way. Write $`\widehat{V}`$ or $`VV`$ for this module.
0.2 Theorem. Suppose the action $`(Q:V)`$ satisfies the following conditions:
(1) $`\text{𝕜}[V]^Q`$ is a polynomial algebra;
(2) $`\mathrm{max}dim_{vV}Qv=dimVdimV//Q`$;
(3) If $`\pi _Q:VV//Q`$ is the quotient morphism and $`\mathrm{\Omega }:=\{vV(\text{d}\pi _Q)_v\text{ is onto }\}`$, then $`V\mathrm{\Omega }`$ contains no divisors.
Set $`\widehat{q}=qq`$ and $`\widehat{Q}=Qq`$. Then
* $`\text{𝕜}[\widehat{V}]^{1q}`$ is a polynomial algebra of Krull dimension $`dimV+dimV//Q`$, which is generated by the coordinates on the first factor of $`\widehat{V}`$ and the polynomials $`\widehat{F}_1,\mathrm{},\widehat{F}_m`$ associated with the differentials of basic invariants in $`\text{𝕜}[V]^Q`$;
* $`\text{𝕜}[\widehat{V}]^{\widehat{Q}}`$ is a polynomial algebra of Krull dimension 2$`dimV//Q`$, which is freely generated by the basic invariants of $`\text{𝕜}[V]^Q`$ and the same functions $`\widehat{F}_i`$, $`i=1,\mathrm{},m`$.
* The $`\widehat{Q}`$-module $`\widehat{V}`$ satisfies conditions (1)–(3), too.
Since the adjoint representation of a reductive Lie algebra $`g`$ satisfies the above properties (1)-(3), one may begin with $`q=g=V`$, and iterate the procedure ad infinitum. For the adjoint representation of a semisimple Lie algebra, the assertion in part (ii) is due to Takiff . For this reason Lie algebras of the form $`qq`$ are called Takiff (Lie) algebras. We will also say that the $`\widehat{q}`$-module $`\widehat{V}`$ is the Takiffisation of the $`q`$-module $`V`$. But $`(g,\mathrm{ad})`$ is not the only possible point of departure for the infinite iteration process. In view of Theorem Introduction, the algebras $`q=gV`$ and their adjoint representations can also be used as initial bricks in the Takiffisation procedure.
If $`\text{𝕜}[V]^Q`$ is polynomial, then it is natural to study the fibres of the quotient morphism $`\pi _Q`$. The null-cone, $`N(V)=\pi _Q^1(\pi _Q(0))`$, is the most important fibre. For instance, $`\text{𝕜}[V]`$ is a free $`\text{𝕜}[V]^Q`$-module if and only $`dimN(V)=dimVdimV//Q`$, i.e., $`\pi _Q`$ is equidimensional. We consider properties of null-cones arising in the context of semi-direct products and their representations.
For $`q=gV`$, as in Theorem Introduction, a necessary and sufficient condition for the equidimensionality of $`\pi _Q`$ is stated in terms of a stratification of $`𝒩`$ determined by the covariants on $`g`$ of type $`V^{}`$. Using this stratification and some technique from and , we prove the following:
If $`N(q)`$ is irreducible, then (i) $`\pi _Q`$ is equidimensional; (ii) the morphism $`\kappa :qq`$ defined by $`\kappa (x,v)=(x,xv)`$, $`xg,vV`$, has the property that the closure of $`\mathrm{Im}(\kappa )`$ is a factorial complete intersection and its ideal in $`\text{𝕜}[q]`$ is generated by the polynomials $`\widehat{F}_i`$, $`i=1,\mathrm{},m`$, mentioned in Theorem Introduction. This is a generalisation of \[22, Prop. 2.4\]. Similar results hold for the Takiffisation of $`G`$-modules $`V`$ having good properties, as in Theorem Introduction. In this case, conditions of equidimensionality for $`\pi _{\widehat{G}}:\widehat{V}\widehat{V}//\widehat{G}`$ are stated in terms of a stratification of $`N(V)`$ determined by the covariants on $`V`$ of type $`V^{}`$. See Section 8 for the details.
In general, it is difficult to deal with the stratifications of $`𝒩`$ and $`N(V)`$, but, for isotropy contractions and $`_2`$-contractions of reductive Lie algebras, explicit results can be obtained. Let $`h`$ be a reductive subalgebra of $`g`$ and $`g=hm`$ a direct sum of $`h`$-modules. Then $`hm`$ is called an isotropy contraction of $`g`$. If $`g=hm`$ is a $`_2`$-grading, then we say about a $`_2`$-contraction. (The word “contraction” can be understood in the usual sense of deformation theory of Lie algebras.) Semi-direct products occurring in this way have some interesting properties. As a sample, we mention the following useful fact: $`\mathrm{ind}(hm)=\mathrm{ind}g+2c(G/H)`$, where $`\mathrm{ind}(.)`$ is the index of a Lie algebra and $`c(.)`$ is the complexity of a homogeneous space. In particular, $`\mathrm{ind}(hm)=\mathrm{ind}g`$ if and only if $`H`$ is a spherical subgroup of $`G`$.
Our main results on the equidimensionality of quotient morphisms and irreducibility of null-cones are related to the $`_2`$-contractions of simple Lie algebras. Given a $`_2`$-grading $`g=g_0g_1`$, Theorem Introduction applies to the semi-direct product $`k=g_0g_1`$, so that $`\text{𝕜}[k]^K`$ is a polynomial algebra of Krull dimension $`\mathrm{rk}g`$. Using the classification of $`_2`$-gradings, we prove that $`N(k)`$ is irreducible. Therefore the good properties discussed in a preceding paragraph hold for the morphism $`\kappa :kk`$, $`\kappa (x_0,x_1)=(x_0,[x_0,x_1])`$. Our proof of irreducibility of $`N(k)`$ basically reduces to the verification of certain inequality for the nilpotent $`G_0`$-orbits in $`g_0`$. Actually, we notice that one may prove a stronger constraint (cf. inequalities (9.8) and (9.9)). This leads to the following curious result: Consider $`\stackrel{~}{k}=g_0(g_1g_1)`$. (In view of Theorem Introduction, $`\text{𝕜}[\stackrel{~}{k}]^{\stackrel{~}{K}}`$ is polynomial.) Then $`\pi _{\stackrel{~}{K}}`$ is still equidimensional, although $`N(\stackrel{~}{k})`$ can already be reducible.
To discuss similar results for the Takiffisation of $`q`$-modules, i.e., $`\widehat{q}`$-modules $`\widehat{V}`$, one has to impose more constraints on $`V`$. We also assume below that $`q=g`$ is reductive.
0.3 Theorem. Suppose the $`G`$-module $`V`$ satisfies conditions (1)–(3) of Theorem Introduction and also the following two conditions:
(4) $`N(V):=\pi _G^1(\pi _G(0))`$ consists of finitely many $`G`$-orbits;
(5) $`N(V)`$ is irreducible and has only rational singularities.
For $`\pi _{\widehat{G}}:\widehat{V}\widehat{V}//\widehat{G}`$ and $`N(\widehat{V})=\pi _{\widehat{G}}^1(\pi _{\widehat{G}}(0))`$, we then have, in addition to the conclusions of Theorem Introduction,
* $`N(\widehat{V})`$ is an irreducible complete intersection and the ideal of $`N(\widehat{V})`$ in $`\text{𝕜}[\widehat{V}]`$ is generated by the basic invariants in $`\text{𝕜}[\widehat{V}]^{\widehat{G}}`$;
* $`\pi _{\widehat{G}}`$ is equidimensional and $`\text{𝕜}[\widehat{V}]`$ is a free $`\text{𝕜}[\widehat{V}]^{\widehat{G}}`$-module.
For $`G`$ semisimple, conditions (2) and (3) are satisfied for all $`V`$, therefore the most essential conditions are (4) and (5). The main point here is to prove the irreducibility. The crucial step in proving this theorem is the use of the Goto-Watanabe inequality \[25, Theorem 2’\] which relates the dimension and embedding dimension of the local rings that are complete intersections with only rational singularities, see Section 10. (We refer to for the definition of rational singularities.) For $`V=g`$, the idea of using that inequality is due to M. Brion. The irreducibility of $`N(\widehat{g})`$ was first proved by F. Geoffriau via case-by-case checking. Then, applying the Goto-Watanabe inequality, Brion found a conceptual proof of Geoffriau’s result . Our observation is that Brion’s idea applies in a slightly more general setting of the Takiffisation of representations $`(G:V)`$ satisfying conditions (1)–(5).
The irreducibility of $`N(\widehat{g})`$ is equivalent to that a certain inequality holds for all non-regular nilpotent elements (orbits). Here is it:
$$dimz_g(x)+\mathrm{rk}(\text{d}\pi _G)_x>2\mathrm{r}\mathrm{k}g\text{ if }x𝒩𝒩^{reg}.$$
Using case-by-case checking, we prove a stronger inequality
$$dimz_g(x)+2\mathrm{r}\mathrm{k}(\text{d}\pi _G)_x3\mathrm{r}\mathrm{k}g0\text{ for all }x𝒩.$$
It seems that the last inequality is more fundamental, because it is stated more uniformly, can be written in different equivalent forms, and has geometric applications. For instance, if $`g=g_0g_1`$ is a $`_2`$-grading of maximal rank and $`\widehat{g}_1=g_1g_1`$, then the equidimensionality of $`\pi _{\widehat{G}_0}:\widehat{g}_1\widehat{g}_1//\widehat{G}_0`$ is essentially equivalent to the last inequality. This result cannot be deduced from Theorem Introduction, because $`N(g_1)`$ is not normal. Furthermore, $`N(\widehat{g}_1)`$ can be reducible.
Our methods also work for generalised Takiff algebras introduced in . The vector space $`q_{\mathrm{}}:=q\text{𝕜}[𝖳]`$ has a natural Lie algebra structure such that $`[x𝖳^l,y𝖳^k]=[x,y]𝖳^{l+k}`$. Then $`q_{(n+1)}={\displaystyle \underset{jn+1}{}}q𝖳^j`$ is an ideal of $`q_{\mathrm{}}`$, and the respective quotient is a generalised Takiff Lie algebra, denoted $`qn`$. Write $`Qn`$ for the corresponding connected group. Clearly, $`dimqn=(n+1)dimq`$ and $`q1qq`$. We prove that if $`(Q:q)`$ satisfies conditions (1)–(3) of Theorem Introduction, then the similar conclusions hold for the adjoint action $`(Qn:qn)`$. In particular, $`\text{𝕜}[qn]^{Qn}`$ is a polynomial algebra of Krull dimension $`(n+1)dimq//Q`$.
For $`q=g`$ semisimple, our methods enable us to deduce the equidimensionality of $`\pi _{G2}:g2g2//G2`$ from the same fact related to the semi-direct product $`g(gg)`$. However, it was shown by Eisenbud and Frenkel that $`\pi _{Gn}:gngn//Gn`$ is equidimensional for any $`n`$, see \[24, Appendix\]. Their proof exploits the interpretation of $`N(gn)`$ as a jet scheme and uses the deep result of Mustaţă concerning the irreducibility of jet schemes \[24, Theorem 3.3\].
Acknowledgements. Work on this article commenced during my visits to the Université de Poitiers (France) in 1996–98. I would like to thank Thierry Levasseur for arranging those visits, inspiring conversations, and drawing my attention to work of Geoffriau. Thanks are also due to Michel Brion for sharing some important insights on Takiff algebras. I am grateful to Sasha Premet for drawing my attention to results of Eisenbud and Frenkel.
## 1. Preliminaries
Algebraic groups are denoted by capital Latin letters and their Lie algebras are denoted by the corresponding lower-case Gothic letters. The identity component of an algebraic group $`Q`$ is denoted by $`Q^o`$.
Let $`Q`$ be an affine algebraic group acting regularly on an irreducible variety $`X`$. Then $`Q_x`$ stands for the isotropy group of $`xX`$. Likewise, the stabiliser of $`x`$ in $`q=\mathrm{Lie}Q`$ is denoted by $`q_x`$. We write $`\text{𝕜}[X]^Q`$ (resp. $`\text{𝕜}(X)^Q`$) for the algebra of regular (resp. field of rational) $`Q`$-invariants on $`X`$. A celebrated theorem of M. Rosenlicht says that there is a dense open $`Q`$-stable subset $`\stackrel{~}{\mathrm{\Omega }}X`$ such that $`\text{𝕜}(X)^H`$ separates the $`Q`$-orbits in $`\stackrel{~}{\mathrm{\Omega }}`$, see e.g. \[5, 1.6\], \[46, 2.3\]. In particular, $`\mathrm{trdeg}\text{𝕜}(X)^Q=dimX\mathrm{max}dim_{xX}Qx`$. We will use Rosenlicht’s theorem in the following equivalent form:
1.1 Theorem. Let $`𝔽`$ be a subfield of $`\text{𝕜}(X)^Q`$. Then $`𝔽=\text{𝕜}(X)^Q`$ if and only if $`𝔽`$ separates the $`Q`$-orbits in a dense open subset of $`X`$. We say that the action $`(Q:X)`$ has a generic stabiliser, if there exists a dense open subset $`\mathrm{\Omega }X`$ such that all stabilisers $`q_\xi `$, $`\xi \mathrm{\Omega }`$, are $`Q`$-conjugate. Then each of the subalgebras $`q_\xi `$, $`\xi \mathrm{\Omega }`$, is called a generic stabiliser. The points of such an $`\mathrm{\Omega }`$ are said to be generic. Likewise, one defines a generic isotropy group, which is a subgroup of $`Q`$. Clearly, the existence of a generic isotropy group implies that of a generic stabiliser. That the converse is also true is proved by Richardson \[34, § 4\]. The reader is also referred to \[46, §7\] for a thorough discussion of generic stabilisers. If $`YX`$ is irreducible, then $`Y^{reg}:=\{yYdimQy=\mathrm{max}_{zY}dimQz\}`$. It is a dense open subset of $`Y`$. The points of $`Y^{reg}`$ are said to be regular. Of course, these notions depend on $`q`$. If we wish to make this dependence explicit, we speak about $`q`$-generic or $`q`$-regular points. Since $`X^{reg}`$ is dense in $`X`$, all generic points (if they do exist) are regular. The converse is however not true.
If $`Q`$ is reductive and $`X`$ is smooth, then $`(Q:X)`$ always has a generic stabiliser . One of our goals is to study existence of generic stabilisers in case of non-reductive $`Q`$. Specifically, we consider the adjoint and coadjoint representations of $`Q`$. To this end, we recall some standard invariant-theoretic techniques and a criterion for the existence of generic stabilisers.
Let $`\rho :QGL(V)`$ be a finite-dimensional rational representation of $`Q`$ and $`\overline{\rho }:qgl(V)`$ the corresponding representation of $`q`$. For $`sQ`$ and $`vV`$, we usually write $`sv`$ in place of $`\rho (s)v`$. Similarly, $`xv`$ is a substitute for $`\overline{\rho }(x)v`$, $`xq`$. (But for the adjoint representation, the standard bracket notation is used.) It should be clear from the context which meaning of ‘$``$’ is meant. Given $`vV`$, consider
$$U=V^{q_v}=\{yVq_vy=0\},$$
the fixed point space of $`q_v`$. Associated to $`UV`$, there are two subgroups of $`Q`$:
$$𝖭(U)=\{sQsUU\},𝖹(U)=\{sQsu=u\text{ for all }uU\}.$$
The following is well known and easy. 1.2 Lemma.
* $`\mathrm{Lie}𝖹(U)=q_v`$ and $`𝖹(U)`$ is a normal subgroup of $`𝖭(U)`$;
* $`𝖭(U)=N_Q(𝖹(U))=N_Q(q_v)`$.
It is not necessarily the case that $`𝖹(U)`$ is connected; however, $`𝖹(U)`$ and $`𝖹(U)^o`$ have the same normaliser in $`Q`$.
1.3 Lemma. If $`yU^{reg}`$ (i.e., $`q_y=q_v`$), then $`QyU=𝖭(U)y`$ and $`qyU=n_q(q_v)y`$. Proof. 1. Suppose $`syU`$ for some $`sQ`$. Then $`q_{sy}=q_v=q_y`$. Hence $`sN_Q(q_v)`$, and we refer to Lemma 1.
2. Suppose $`syU`$ for some $`sq`$. Then $`0=q_v(sy)=[q_v,s]y`$. Hence $`[q_v,s]q_y=q_v`$. Set $`Y=\overline{QU}`$. It is a $`Q`$-stable irreducible subvariety of $`V`$.
1.4 Proposition. The restriction homomorphism $`(f\text{𝕜}(Y))f|_U`$ yields an isomorphism $`\text{𝕜}(Y)^Q\stackrel{}{}\text{𝕜}(U)^{𝖭(U)}=\text{𝕜}(U)^{𝖭(U)/𝖹(U)}`$. Proof. This follows from the first equality in Lemma 1 and Rosenlicht’s theorem. 1.5 Example. Let $`G`$ be a semisimple algebraic group with Lie algebra $`g`$, and $`v=eg`$ a nilpotent element. Then $`g_e=z_g(e)`$ is the centraliser of $`e`$ and $`U=\{xg[x,z_g(e)]=0\}=:d_g(e)`$ is the centre of $`z_g(e)`$. Here $`𝖭(U)=N_G(z_g(e))`$ is the normaliser of $`z_g(e)`$ in $`G`$. Letting $`Y=\overline{Gd_g(e)}`$, we obtain an isomorphism
$$\text{𝕜}(Y)^G\text{𝕜}(d_g(e))^{N_G(z_g(e))}.$$
It is known that $`d_g(e)`$ contains no semisimple elements , so that $`Y`$ is the closure of a nilpotent orbit and hence $`\text{𝕜}(Y)^G=\text{𝕜}`$. It follows that $`N_G(z_g(e))`$ has a dense orbit in $`d_g(e)`$. This fact was already noticed in \[30, § 4\]. Actually, the dense $`G`$-orbit in $`Y`$ is just $`Ge`$. Clearly, if $`\overline{QU}=V`$, then $`(Q:V)`$ has a generic stabiliser and $`v`$ is a generic point. A general criterion for this to happen is proved in \[14, § 1\]. For future reference, we recall it here.
1.6 Lemma (Elashvili). Let $`vV`$ be an arbitrary point. Then $`QV^{q_v}`$ is dense in $`V`$ if and only if $`V=qv+V^{q_v}`$. The existence of a non-trivial generic stabiliser yields a Chevalley-type theorem for the field of invariants. Indeed, it follows from Proposition 1 that if $`(Q:V)`$ has a generic stabiliser, $`vV`$ is a generic point, and $`U=V^{q_v}`$, then
(1.7)
$$\text{𝕜}(V)^Q\text{𝕜}(U)^{𝖭(U)}=\text{𝕜}(U)^{𝖭(U)/𝖹(U)}.$$
In this context, the group $`W:=𝖭(U)/𝖹(U)`$ is called the Weyl group of the action $`(Q:V)`$. Notice that this $`W`$ is not necessarily finite.
The corresponding question for the algebras of invariants is much more subtle. The restriction homomorphism $`ff|_U`$ certainly induces an embedding $`\text{𝕜}[V]^Q\text{𝕜}[U]^{𝖭(U)/𝖹(U)}`$. However, if $`Q`$ is non-reductive, then it is usually not onto.
## 2. Generic stabilisers (centralisers) for the adjoint representation
In what follows, $`Q`$ is a connected algebraic group. In this section, we elaborate on the existence of generic stabilisers and its consequences for the adjoint representations $`\mathrm{Ad}:QGL(q)`$ and $`\mathrm{ad}:q𝔤𝔩(q)`$.
For $`xq`$, the stabiliser $`q_x`$ is nothing but the centraliser of $`x`$ in $`q`$, so that we write $`z_q(x)`$ in place of $`q_x`$. The centraliser of $`x`$ in $`Q`$ is denoted by $`Z_Q(x)`$. If $`(q,\mathrm{ad})`$ has a generic stabiliser, then we also say that $`q`$ has a generic centraliser. By Lemma 1, a point $`xq`$ is generic if and only if
$$[q,x]+q^{z_q(x)}=q.$$
Since $`q^{z_q(x)}`$ is the centre of the Lie algebra $`z_q(x)`$ and $`dim[q,x]=dimqdimz_q(x)`$, one immediately derives
2.1 Proposition. An algebraic Lie algebra $`q`$ has a generic centraliser if and only if there is an $`xq`$ such that
(2.2) $`z_q(x)`$ is commutative and
(2.3) $`[q,x]z_q(x)=q`$.
Equality (2.3) implies that $`\mathrm{Im}(\mathrm{ad}x)=\mathrm{Im}(\mathrm{ad}x)^2`$. The latter is never satisfied if $`\mathrm{ad}x`$ is nilpotent and $`\mathrm{Im}(\mathrm{ad}x)0`$. That is, if $`q`$ is nilpotent and $`[q,q]0`$, then $`q`$ has no generic centralisers. It also may happen that neither of the centralisers $`z_q(x)`$ is commutative. (Consider the Heisenberg Lie algebra $`H_n`$ of dimension $`2n+1`$ for $`n2`$.) On the other hand, if there is a semisimple $`xq`$ such that $`z_q(x)`$ is commutative, then the conditions of Proposition 2 are satisfied, so that a generic centraliser exists. \[Warning: this does not imply that the semisimple elements are dense in $`q`$.\] 2.4 Lemma. Let $`xq`$ be a generic point. Then $`n_q(z_q(x))=z_q(x)`$. Proof. Assume that $`n_q(z_q(x))z_q(x)`$. In view of Eq. (2.3), there is then a nonzero $`yn_q(z_q(x))[q,x]`$. That is, $`y=[s,x]`$ for some $`sq`$. Then
$$[y,z_q(x)]=[[s,z_q(x)],x][q,x]$$
and hence $`[y,z_q(x)]=0`$. Thus, $`yz_q(x)[q,x]=0`$, and we are done. Recall that a subalgebra $`h`$ of $`q`$ is called a Cartan subalgebra if $`h`$ is nilpotent and $`n_q(h)=h`$. Every Lie algebra has a Cartan subalgebra, and all Cartan subalgebras of $`q`$ are conjugate under $`Q`$, see \[37, Ch. III\].
2.5 Proposition. An algebraic Lie algebra $`q`$ has a generic centraliser if and only if the Cartan subalgebras of $`q`$ are commutative. Proof. If $`q`$ has a generic centraliser, then, by Lemma 2, such a centraliser is a (commutative) Cartan subalgebra. Conversely, any Cartan subalgebra of $`q`$ is of the form $`h=\{yq(\mathrm{ad}x)^ny=0\text{ for }n0\}`$ for some $`xq`$ \[37, Ch. III.4, Cor. 2\]. Therefore, the commutativity of $`h`$ implies that $`h=z_q(x)`$ and $`\mathrm{ad}x`$ is invertible on $`[q,x]`$. As is already mentioned, the existence of a generic centraliser implies that of a generic isotropy group. For this reason, we always assume that a generic point $`x`$ has the property that $`Z_Q(x)`$ is a generic isotropy group. (This is only needed if a generic isotropy group is disconnected.)
2.6 Theorem. Suppose $`q`$ has a generic centraliser. Let $`xq`$ be a generic point such that $`Z_Q(x)`$ is a generic isotropy group. Then (i) $`𝖹(z_q(x))=Z_Q(x)`$ and (ii) $`\text{𝕜}(q)^Q\text{𝕜}(z_q(x))^W`$, where $`W=N_Q(z_q(x))/Z_Q(x)`$ is a finite group. Proof. (i) Since $`xz_q(x)`$, we have $`𝖹(z_q(x))Z_Q(x)`$. Hence one has to prove that $`Z_Q(x)`$ acts trivially on $`z_q(x)`$. Assume that the fixed point space of $`Z_Q(x)`$ is a proper subspace of $`z_q(x)`$, say $`M`$. Since $`dimQMdim[q,x]+dimM<dimq`$, $`QM`$ cannot be dense in $`q`$, which contradicts the fact that $`Z_Q(x)`$ is a generic isotropy group.
(ii) This follows from Eq. (1.7) and Lemma 2. Below, we state a property of generic points related to the dual space $`q^{}`$.
2.7 Proposition. Let $`xq`$ be a generic point, as in Theorem 2. Then
(i) $`q^{}=xq^{}(q^{})^x=xq^{}(q^{})^{z_q(x)}`$ and (ii) $`(q^{})^{Z_Q(x)}=(q^{})^{z_q(x)}`$. Proof. (i) We have $`[q,x]^{}=(q^{})^x`$ and $`z_q(x)^{}=xq^{}`$. Hence the first equality follows from Eq. (2.3).
The second equality means that $`(q^{})^x=(q^{})^{z_q(x)}`$. Clearly, $`(q^{})^x(q^{})^{z_q(x)}`$. Taking the annihilators provides the inclusion $`[q,x][q,z_q(x)]`$. Then using Eq. (2.2) and (2.3) yields
$`[q,z_q(x)][z_q(x)+[q,x],z_q(x)]=[[q,x],z_q(x)]=[[q,z_q(x)],x][q,x]`$.
(ii) In view of (i), $`(q^{})^{z_q(x)}`$ is identified with $`(z_q(x))^{}`$. Hence the assertion stems from Theorem 2(i). Thus, the very existence of a generic centraliser implies that $`q`$ has some properties in common with reductive Lie algebras. For instance, the Weyl group of $`(Q:q)`$ is finite, and the decomposition of $`q^{}`$ with respect to a generic element $`xq`$ is very similar to that of $`q`$. It will be shown below that there is a vast stock of such Lie algebras.
## 3. Generic stabilisers for the coadjoint representation
In this section, we work with the coadjoint representations of $`Q`$ and $`q`$. Usually, we use lowercase Latin (resp. Greek) letters to denote elements of $`q`$ (resp. $`q^{}`$). By Lemma 1, a point $`\xi q^{}`$ is generic if and only if
$$q\xi +(q^{})^{q_\xi }=q^{}.$$
As was noticed by Tauvel and Yu , taking the annihilators yields a simple condition, entirely in terms of $`q`$. Namely, $`\xi `$ is generic if and only if
(3.1)
$$q_\xi [q,q_\xi ]=\{0\}.$$
Below, we assume that $`(q,\mathrm{ad}^{})`$ has a generic stabiliser and thereby Eq. (3.1) is satisfied for some $`\xi `$. This readily implies that $`q_\xi `$ is commutative and $`n_q(q_\xi )=z_q(q_\xi )`$. However, unlike the adjoint representation case, $`q_\xi `$ can be a proper subalgebra of $`z_q(q_\xi )`$. In other words, the Weyl group of $`(Q:q^{})`$ is not necessarily finite. Our goal is to understand what isomorphism (1.7) means in this situation. Set $`h=q_\xi `$ and $`U=(q^{})^{q_\xi }`$. Then we can write
$$\text{𝕜}(q^{})^Q(\text{𝕜}(U)^{Z_Q(h)^o})^{N_Q(h)/Z_Q(h)^o}.$$
That is, one first takes the invariants of the connected group $`Z_Q(h)^o`$, and then the invariants of the finite group $`N_Q(h)/Z_Q(h)^o`$. 3.2 Lemma. $`dimU=dimz_q(h)`$. Proof. By Lemma 1 and Eq. (3.1), we have $`q\xi U=z_q(h)\xi `$. Equating the dimensions of these spaces yields the assertion. In view of this equality, it is tempting to interpret $`U`$ as the space of the coadjoint representation of $`z_q(h)=\mathrm{Lie}Z_Q(h)^o`$. However it seems to only be possible under an additional assumption on $`h`$.
3.3 Definition. We say that a subalgebra $`h`$ is near-toral if $`[q,h]z_q(h)=\{0\}`$. This condition is stronger than (3.1). It is obviously satisfied if $`h`$ is a toral Lie algebra (= Lie algebra of a torus).
Recall that the index of (a Lie algebra) $`q`$, $`\mathrm{ind}q`$, is the minimal codimension of $`Q`$-orbits in $`q^{}`$. Equivalently, $`\mathrm{ind}q=\mathrm{trdeg}\text{𝕜}(q^{})^Q`$. If $`\mathrm{ind}q=0`$, then $`q`$ is called Frobenius.
3.4 Theorem. Suppose the generic stabiliser $`h`$ is near-toral. Then
* $`[q,h]z_q(h)=q`$ and $`Uz_q(h)^{}`$;
* $`\mathrm{ind}q=\mathrm{ind}z_q(h)=dimh`$ and $`h`$ is the centre of $`z_q(h)`$
Proof. (i) It is easily seen that $`[q,h]^{}=(q^{})^h=U`$. Therefore Definition 3 says that $`z_q(h)^{}+U=q^{}`$. From Lemma 3, it then follows that this sum (of $`z_q(h)`$-modules) is direct. Hence $`Uq^{}/z_q(h)^{}z_q(h)^{}`$.
(ii) Since $`\xi `$ is generic and hence regular in $`q^{}`$, we have $`\mathrm{ind}q=dimh`$.
For $`\nu U^{reg}`$, we have $`Uh^{}=Uq\nu =z_q(h)\nu `$. In particular, $`dimz_q(h)\nu =dimUdimh`$. Hence almost all $`Z_Q(h)`$-orbits in $`U`$ are of codimension $`dimh`$. This also means that the centre of $`z_q(h)`$ cannot be larger than $`h`$. 3.5 Corollary. If the generic stabiliser $`h`$ is near-toral, then $`\text{𝕜}(q^{})^Q(\text{𝕜}(z_q(h)^{})^{Z_Q(h)^o})^F`$, where $`F=N_Q(h)/Z_Q(h)^o`$ is finite. That is, one first takes the invariants of the coadjoint representation for a smaller Lie algebra and then the invariants of a finite group. Under the assumption that $`h`$ is near-toral, $`s:=z_q(h)`$ has the property that $`\mathrm{ind}s=dimz(s)`$. The following results present some properties of such algebras.
3.6 Proposition. Suppose $`\mathrm{ind}s=dimz(s)`$. Then
1. The closure of any regular $`S`$-orbit in $`s^{}`$ is an affine space.
2. If $`z(s)`$ is toral, then $`s/z(s)`$ is Frobenius. Proof. 1. If $`y(s^{})^{reg}`$, then $`s_y=z(s)`$ and hence $`sy=z(s)^{}`$. Hence all points of the orbit $`Sy`$ have one and the same tangent space. Therefore $`Sy`$ is open and dense in the affine space $`y+z(s)^{}`$.
2. Since $`z(s)`$ is reductive, one has a direct sum of Lie algebras $`s=rz(s)`$, and $`\mathrm{ind}r=\mathrm{ind}s\mathrm{ind}z(s)=0`$. It is not, however, always true that $`s/z(s)`$ is Frobenius. For instance, the Heisenberg Lie algebra $`H_n`$ has one-dimensional centre and $`\mathrm{ind}H_n=1`$. But $`H_n/z(H_n)`$ is commutative, so that $`\mathrm{ind}(H_n/z(H_n))=2n`$.
3.7 Examples. 1. Let $`b`$ be a Borel subalgebra of a simple Lie algebra $`g`$. Then $`(b,\mathrm{ad}^{})`$ has a generic stabiliser, which is always a toral Lie algebra, see e.g. . If $`h`$ is such a stabiliser, then by Proposition 3, $`z_b(h)/h`$ is a Frobenius Lie algebra. It is not hard to compute this quotient for all cases in which $`h0`$.
$``$ If $`g=sl_n`$, then $`dimh=\left[\frac{n1}{2}\right]`$ and $`z_b(h)/hb(sl_2)^{[n/2]}`$.
$``$ If $`g=so_{4n+2}`$, then then $`dimh=1`$ and $`z_b(h)/hb(so_{4n})`$.
$``$ If $`g=\text{E}_6`$, then $`dimh=2`$ and $`z_b(h)/hb(so_8)`$.
2. If $`g=sl_n`$ or $`sp_{2n}`$ and $`s`$ is a seaweed subalgebra of $`g`$, then a generic stabiliser for $`(s,\mathrm{ad}^{})`$ always exists, and it is a toral subalgebra . For instance, let $`pgl_{2n}`$ be a maximal parabolic subalgebra whose Levi part is $`gl_ngl_n`$. Then a generic stabiliser for $`(p,\mathrm{ad}^{})`$ is $`n`$-dimensional and toral, and $`z_p(h)/hb(sl_2)^n`$.
3. There are non-trivial examples of Lie algebras such that a generic stabiliser for $`\mathrm{ad}^{}`$ exists, is near-toral, and equals its own centraliser, but it is not toral. Let $`e`$ be a nilpotent element in $`g=sl_n`$ and $`q=z_g(e)`$. Then a generic stabiliser for the coadjoint representation of $`q`$ exists, see . If $`h`$ is such a stabiliser, then the description of $`h`$ given in \[48, Theorems 1 & 5\] shows that $`z_q(h)=h`$. Hence, by Corollary 3, $`\text{𝕜}(q^{})^Q`$ is the field of invariants of a finite group.
## 4. Semi-direct products of Lie algebras and modules of covariants
In this section, we review some notions and results that will play the principal role in the following exposition.
(I) Recall a semi-direct product construction for Lie groups and algebras.
Let $`V`$ be a $`Q`$-module, and hence a $`q`$-module. Then $`q\times V`$ has a natural structure of Lie algebra, $`V`$ being an Abelian ideal in it. Explicitly, if $`x,x^{}q`$ and $`v,v^{}V`$, then
$$[(x,v),(x^{},v^{})]=([x,x^{}],xv^{}x^{}v).$$
This Lie algebra is denoted by $`qV`$ or $`qϵV`$. Accordingly, an element of this algebra is denoted by either $`(x,v)`$ or $`x+ϵv`$. Here $`ϵ`$ is regarded as a formal symbol. Sometimes, e.g. if $`V=q`$, it is convenient to think of $`ϵ`$ as element of the ring of dual numbers $`\text{𝕜}[ϵ]=\text{𝕜}\text{𝕜}ϵ`$, $`ϵ^2=0`$. A connected algebraic group with Lie algebra $`qV`$ is identified set-theoretically with $`Q\times V`$, and we write $`QV`$ for it. The product in $`QV`$ is given by
$$(s,v)(s^{},v^{})=(ss^{},(s^{})^1v+v^{}).$$
In particular, $`(s,v)^1=(s^1,sv)`$. The adjoint representation of $`QV`$ is given by the formula
(4.1)
$$(\mathrm{Ad}(s,v))(x^{},v^{})=(\mathrm{Ad}(s)x^{},sv^{}x^{}v),$$
where $`v,v^{}V`$, $`xq`$, and $`sQ`$.
Note that $`V`$ can be regarded as either a commutative unipotent subgroup of $`QV`$ or a commutative nilpotent subalgebra of $`qV`$. Referring to $`V`$ as subgroup of $`QV`$, we write $`1V`$. A semi-direct product $`qV`$ is said to be reductive if $`q`$ is a reductive (algebraic) Lie algebra.
(II) Our second important ingredient is the notion of modules of covariants.
Let $`A`$ be an algebraic group, acting on an affine variety $`X`$, and $`V`$ an $`A`$-module. The set of all $`A`$-equivariant morphisms from $`X`$ to $`V`$, denoted $`\mathrm{Mor}_A(X,V)`$, has a natural structure of $`\text{𝕜}[X]^A`$-module. This $`\text{𝕜}[X]^A`$-module is said to be the module of covariants (of type $`V`$). It is easily seen that $`\mathrm{Mor}_A(X,V)`$ can be identified with $`(\text{𝕜}[X]V)^A`$. For any $`xX`$, we denote by $`\epsilon _x`$ the evaluation homomorphism $`\mathrm{Mor}_A(X,V)V`$, which takes $`F`$ to $`F(x)`$. Obviously, $`\mathrm{Im}(\epsilon _x)V^{A_x}`$.
Assume for a while that $`A=G`$ is reductive. Then the algebra $`\text{𝕜}[X]^G`$ is finitely generated and $`\mathrm{Mor}_G(X,V)`$ is a finitely generated $`\text{𝕜}[X]^G`$-module, see e.g. \[5, 2.5\], \[46, 3.12\]. A review of recent results on modules of covariants in the reductive case can be found in . The following result is proved in \[29, Theorem 1\].
4.2 Theorem. If $`\overline{Gx}`$ is normal and $`\mathrm{codim}_{\overline{Gx}}(\overline{Gx}Gx)2`$, then $`\mathrm{Im}(\epsilon _x)=V^{G_x}`$. Let $`g^{reg}`$ be the set of regular elements of $`g`$ and $`T`$ a maximal torus of $`G`$. The following fundamental result is due to Kostant \[21, p. 385\].
4.3 Theorem. Let $`V`$ be a $`G`$-module. Then $`dimV^{G_x}=dimV^T`$ for any $`xg^{reg}`$ and $`\mathrm{Mor}_G(g,V)`$ is a free $`\text{𝕜}[g]^G`$-module of rank $`dimV^T`$. In particular, if $`V^T=0`$, then there is no non-trivial $`G`$-equivariant mappings from $`g`$ to $`V`$. These modules of covariants are graded, and the degrees of minimal generating systems are uniquely determined. These degrees are called the generalised exponents of $`V`$. The multiset of generalised exponents of a $`g`$-module $`V`$ is denoted by $`\text{g-exp}_g(V)`$. Similar results hold if $`g`$ is replaced with a ”sufficiently good” $`G`$-module, see \[47, Ch. III, § 1\] and \[36, Prop. 4.3, 4.6\]. Namely,
4.4 Theorem. Let $`\stackrel{~}{V}`$ be a $`G`$-module such that $`\text{𝕜}[\stackrel{~}{V}]^G`$ is a polynomial algebra and the quotient morphism $`\pi :\stackrel{~}{V}\stackrel{~}{V}//G`$ is equidimensional. Then $`\mathrm{Mor}_G(\stackrel{~}{V},V)`$ is a free $`\text{𝕜}[\stackrel{~}{V}]^G`$-module for any $`G`$-module $`V`$. Furthermore, if $`(G:\stackrel{~}{V})`$ is stable, then the rank of $`\mathrm{Mor}_G(\stackrel{~}{V},V)`$ equals $`dimV^H`$, where $`H`$ is a generic isotropy group for $`(G:\stackrel{~}{V})`$. An action $`(G:V)`$ is said to be stable, if the union of closed $`G`$-orbits is dense in $`V`$ (see \[46, 7.5\] and about stable actions). If $`(G:V)`$ is stable, then a generic stabiliser is reductive and $`\text{𝕜}(V)^G`$ is the quotient field of $`\text{𝕜}[V]^G`$.
In some cases, a basis for free modules of covariants can explicitly be indicated. For any $`f\text{𝕜}[V]`$, the differential of $`f`$ can be regarded as a covector field on $`V`$: $`v\text{d}f_vV^{}`$. Starting with $`f\text{𝕜}[V]^G`$, one obtains in this way a covariant $`\text{d}f\mathrm{Mor}_G(V,V^{})`$. The following result of Thierry Vust appears in \[47, Ch. III, § 2\].
4.5 Theorem. Let a $`G`$-module $`\stackrel{~}{V}`$ satisfy all the assumptions of Theorem 4. Suppose also that $`N_G(H)/H`$ is finite. Let $`f_1,\mathrm{},f_m`$ be a set of basic invariants in $`\text{𝕜}[\stackrel{~}{V}]^G`$. Then $`\mathrm{Mor}_G(\stackrel{~}{V},\stackrel{~}{V}^{})`$ is freely generated by $`\text{d}f_i`$, $`i=1,\mathrm{},m`$.
(III) Here we point out a connection between modules of covariants and invariants of semi-direct products.
For $`F\mathrm{Mor}_A(X,V)`$, define the polynomial $`\widehat{F}\text{𝕜}[X\times V^{}]^A`$ by the rule $`\widehat{F}(x,\xi )=F(x),\xi `$, where $`,:V\times V^{}\text{𝕜}`$ is the natural pairing.
4.6 Lemma. Consider the Lie algebra $`qV`$ and the $`\text{𝕜}[q]^Q`$-module $`\mathrm{Mor}_Q(q,V^{})`$. Then for any $`F\mathrm{Mor}_Q(q,V^{})`$, we have $`\widehat{F}\text{𝕜}[qV]^{QV}`$. Proof. Clearly, $`\widehat{F}`$ is $`Q`$-invariant. The invariance with respect to $`1V`$-action means that
$$F(x),v=F(x),v+xv^{}$$
holds for any $`xq`$ and $`v,v^{}V`$. To this end, we notice that $`F(x),xv^{}=xF(x),v^{}`$, and $`xF(x)=0`$, since $`F:qV^{}`$ is a $`Q`$-equivariant morphism. The point is that $`\widehat{F}`$ turns out to be invariant with respect to the action of the unipotent group $`1V`$.
## 5. Generic stabilisers and rational invariants for semi-direct products
Given $`Q`$ and $`V`$, one may ask the following questions:
* When does a generic centraliser for $`qV`$ exist? What are invariant-theoretic consequences of this?
It is easily seen that the existence of a generic centraliser for $`q`$ is a necessary condition. We will therefore assume that this is the case.
5.1 Theorem. Let $`xq`$ be a generic point. Suppose $`V^x=V^{Z_Q(x)}`$ and $`V^xxV=V`$. Then
* each point of the form $`x+ϵv`$, $`vV^{z_q(x)}`$, is generic and $`z_q(x)ϵV^{z_q(x)}`$ is a generic centraliser for $`qV`$.
* $`\mathrm{trdeg}\text{𝕜}(qV)^{QV}=\mathrm{trdeg}\text{𝕜}(q)^Q+dimV^{z_q(x)}`$;
* The Weyl groups of $`(q,\mathrm{ad})`$ and $`(qV,\mathrm{ad})`$ are isomorphic;
* $`\text{𝕜}(qV)^{QV}`$ is a purely transcendental extension of $`\text{𝕜}(q)^Q`$.
Proof. Set $`h=z_q(x)`$, $`R=QV`$, and $`r=qV`$. It follows from the assumptions that $`V^x=V^h`$.
(i) Let $`vV^h`$ be arbitrary. Let us verify that Proposition 2 applies here. A direct calculation shows that $`z_r(x+ϵv)=hϵV^h`$ and this algebra is commutative. Next,
$$[r,x+ϵv]=\{[z,x]+ϵ(zv)zq\}+ϵ(xV).$$
Notice that $`qv=([q,x]h)v=[q,x]v=x(qv)xV`$. Hence $`zvxV`$ for any $`zq`$ and $`[r,x+ϵv]=[q,x]ϵ(xV)`$. Therefore the equality $`[r,x+ϵv]z_r(x+ϵv)=r`$ is equivalent to that $`V^xxV=V`$.
(ii) By part (i), $`\stackrel{~}{h}:=hϵV^h`$ is a generic centraliser for $`r`$. Since $`\mathrm{trdeg}\text{𝕜}(q)^Q=dimh`$, the claim follows.
(iii) Using formula (4.1), one easily verifies that $`N_R(\stackrel{~}{h})=N_Q(h)V^h`$ and $`Z_R(\stackrel{~}{h})=Z_Q(h)V^h`$. Hence using Theorem 2, we obtain
$$\stackrel{~}{W}=N_R(\stackrel{~}{h})/Z_R(\stackrel{~}{h})N_Q(h)/Z_Q(h)=W.$$
(iv) Here we may work entirely with invariants of $`W`$. In view of (iii) and Theorem 2, it suffices to prove that $`\text{𝕜}(h)^W\text{𝕜}(hV^h)^W`$ is a purely transcendental extension. Actually, a transcendence basis of $`\text{𝕜}(hV^h)^W`$ over $`\text{𝕜}(h)^W`$ can explicitly be constructed. This follows from Theorem 5 below, since the representation of $`W`$ on $`h`$ is faithful. The following result concerns fields of invariants of reductive algebraic groups.
Recall from Section 4(III) that one may associate the invariant $`\widehat{F}\text{𝕜}[V_1\times V_2]^G`$ to any $`F\mathrm{Mor}_G(V_1,V_2^{})`$. If $`D`$ is a domain, then we write $`D_{(0)}`$ for the field of fractions.
5.2 Theorem. Let $`\rho _i:GGL(V_i)`$, $`i=1,2`$, be representations of a reductive group $`G`$. Set $`m=dimV_2`$ and $`J=\text{𝕜}[V_1]^G`$. Suppose that a generic isotropy subgroup for $`(G:V_1)`$ is trivial, and $`(G:V_1)`$ is stable. Then
* $`dim_{J_{(0)}}\mathrm{Mor}_G(V_1,V_2^{})_JJ_{(0)}=m`$;
* Let $`F_1,\mathrm{},F_m\mathrm{Mor}_G(V_1,V_2^{})`$ be covariants such that $`\{F_i1i=1,\mathrm{},m\}`$ form a basis for the $`J_{(0)}`$-vector space in (i). Then $`\text{𝕜}(V_1V_2)^G=\text{𝕜}(V_1)^G(\widehat{F}_1,\mathrm{},\widehat{F}_m)`$. In other words, any such basis for $`\mathrm{Mor}_G(V_1,V_2^{})_JJ_{(0)}`$ gives rise to a transcendence basis for the field $`\text{𝕜}(V_1V_2)^G`$ over $`\text{𝕜}(V_1)^G`$.
Proof. (i) Because $`(G:V_1)`$ is stable, $`J_{(0)}=\text{𝕜}(V_1)^G`$. Since $`\mathrm{Mor}_G(V_1,V_2^{})`$ is a finitely-generated $`J`$-module, $`\underset{¯}{M}=\mathrm{Mor}_G(V_1,V_2^{})_JJ_{(0)}`$ is a finite-dimensional $`J_{(0)}`$-vector space. By the assumptions, there is an $`xV_1`$ such that the isotropy group $`G_x`$ is trivial and $`Gx=\overline{Gx}`$. Then by Theorem 4,
$`()`$ the evaluation map $`\epsilon _x:\mathrm{Mor}_G(V_1,V_2^{})V_2^{}=(V_2^{})^{G_x}`$ is onto. Hence $`dim\underset{¯}{M}m`$. On the other hand, it cannot be greater than $`m`$.
(ii) In view of Theorem 1, it suffices to prove that $`\text{𝕜}(V_1)^G(\widehat{F}_1,\mathrm{},\widehat{F}_m)`$ separates the generic $`G`$-orbits in $`V_1V_2`$. First, the field $`\text{𝕜}(V_1)^G`$ separates the generic $`G`$-orbits in $`V_1`$. Therefore, for generic points $`(x_1,x_2)`$, $`x_iV_i`$, the first coordinate is determined uniquely up to $`G`$-conjugation by the values $`f(x_1)`$, where $`f`$ runs over $`\text{𝕜}(V_1)^G`$. By condition $`()`$, $`F_1(x_1),\mathrm{},F_m(x_1)`$ form a basis for $`V_2^{}`$ if $`x_1`$ is generic. Hence given a generic $`x_1`$ and arbitrary values of the invariants $`\widehat{F}_i`$, the second coordinate (i.e., $`x_2`$) is uniquely determined. Remarks. 1. Most of the assumptions of Theorem 5 are always satisfied if $`G`$ is either finite or semisimple. For $`G`$ finite, it suffices to only require that $`\rho _1`$ is faithful. For $`G`$ semisimple, it suffices to require that a generic isotropy group of $`(G:V_1)`$ is trivial.
2. The assertion that the field extension in (ii) is purely transcendental is known, see e.g. \[12, p. 6\]. But the explicit construction of a transcendence basis via modules of covariants seems to be new.
The following assertion demonstrates important instances, where Theorem 5 applies.
5.3 Proposition. Theorem 5 applies to the following $`q`$-modules $`V`$:
1. $`q`$ is an arbitrary Lie algebra having a generic centraliser and $`V`$ is either $`q`$ or $`q^{}`$.
2. $`q=g`$ is reductive and $`V`$ is an arbitrary $`g`$-module. Proof. 1. For $`qq`$, the conditions of Theorem 5 are satisfied in view of Proposition 2 and Theorem 2. For $`qq^{}`$, these conditions are satisfied in view of Proposition 2.
2. Here $`xg`$ is a regular semisimple element and $`Z_G(x)`$ is a maximal torus. Therefore $`V^x`$ is the zero weight space of $`V`$ and $`xV`$ is the sum of all other weight spaces. Remark. For the semi-direct products as in Proposition 5(2), we are able to describe the polynomial invariants, see Section 6.
A Lie algebra is said to be quadratic whenever its adjoint and coadjoint representations are equivalent. It is easily seen that $`qq^{}`$ is quadratic for any Lie algebra $`q`$. For, if $`,`$ is the pairing of $`q`$ and $`q^{}`$, then the formula $`(x_1+ϵ\xi _1,x_2+ϵ\xi _2)=x_1,\xi _2+x_2,\xi _1`$ determines a non-degenerate symmetric $`qq^{}`$-invariant form. For $`qq^{}`$, there is no difference between the adjoint and coadjoint representations. So, previous results of this section describe some properties of the coadjoint representation of $`qq^{}`$ as well. However, for an arbitrary $`V`$ the adjoint and coadjoint representation of $`qV`$ are very different. Hence our second problem is:
* When does a generic stabiliser for $`(qV,\mathrm{ad}^{})`$ exist? What are invariant-theoretic consequences of this?
This problem is quite different from (Q1). It seems to be more involved and restrictive.
Set $`r=qV`$ and $`R=QV`$. The dual space $`r^{}`$ is identified with $`q^{}V^{}`$, and a typical element of it is denoted by $`\eta =(\alpha ,\xi )`$. For $`(s,v)r`$, the coadjoint representation is given by
(5.4)
$$(\mathrm{ad}_r^{}(s,v))(\alpha ,\xi )=(\mathrm{ad}_q^{}(s)(\alpha )v\xi ,s\xi ).$$
Here the mapping $`((s,\xi )q\times V^{})(s\xi V^{})`$ is the natural $`q`$-module structure on $`V^{}`$, and $`((v,\xi )V\times V^{})(v\xi q^{})`$ is the moment mapping with respect to the symplectic structure on $`V\times V^{}`$.
To describe the stabiliser of any point in $`r^{}`$, we need some notation. For $`\alpha q^{}`$, let $`𝒦_\alpha `$ denote the Kirillov form on $`q`$ associated with $`\alpha `$, i.e., $`𝒦_\alpha (s_1,s_2)=\alpha ,[s_1,s_2]`$. Then $`\mathrm{ker}(𝒦_\alpha )=q_\alpha `$, the stabiliser of $`\alpha `$. If $`h`$ is a subalgebra of $`q`$, then $`𝒦_\alpha |_h`$ can also be regarded as the Kirillov form associated with $`\alpha |_hh^{}`$.
5.5 Proposition. For any $`\eta =(\alpha ,\xi )r^{}`$, we have
$$r_\eta =\{(s,v)s\mathrm{ker}(𝒦_\alpha |_{q_\xi })\&\mathrm{ad}_q^{}(s)\alpha =v\xi \}.$$
Proof. Straightforward. The first condition imposed on $`s`$ guarantees us the equality $`s\xi =0`$ and that the equation $`\mathrm{ad}_q^{}(s)\alpha =v\xi `$ has a solution $`v`$ for any such $`s`$. It follows that $`r_\eta `$ is a direct sum of the space $`\{wVw\xi =0\}=(q\xi )^{}`$, sitting in $`V`$, and a space of dimension $`dim\mathrm{ker}(𝒦_\alpha |_{q_\xi })`$, which is embedded in $`qV`$ somehow diagonally. (We will see below that under additional constraints this second space lies entirely in $`q`$.)
A result of Raïs on semi-direct products describes $`r`$-regular points in $`r^{}`$ and gives the value of $`\mathrm{ind}r`$, that is, the dimension of the stabilizer of the $`r`$-regular points in $`r^{}`$. Namely, if $`\xi V^{}`$ is $`q`$-regular, then $`(\alpha ,\xi )`$ is $`r`$-regular if and only if $`\alpha `$ is $`q_\xi `$-regular as element of $`q_\xi ^{}`$ (with respect to the coadjoint representation of $`q_\xi `$). By a theorem of Duflo-Vergne , the stabiliser of any regular point in the coadjoint representation is commutative, see also \[30, 1.8\] for an invariant-theoretic proof.
It seems to be difficult to find out a general condition ensuring that Eq. (3.1) holds for some regular point in $`r^{}`$. For this reason, we only look at the three cases occurring already in Proposition 5 in connection with generic centralisers.
$``$ If $`(q,\mathrm{ad}^{})`$ has a generic stabiliser, then $`(qq,\mathrm{ad}^{})`$ has.
Indeed, if $`q_\xi `$ is a generic stabiliser ($`\xi q^{}`$), then $`q_\xi q_\xi `$ is the stabiliser of $`\eta =(0,\xi )(qq)^{}`$ and $`[qq,q_\xi q_\xi ](q_\xi q_\xi )=\{0\}`$.
$``$ If $`(q,\mathrm{ad}^{})`$ has a generic stabiliser, then $`(qq^{},\mathrm{ad}^{})`$ may have no generic stabilisers.
Example. Let $`q`$ be the 3-dimensional Heisenberg algebra $`H_1`$. The generic stabiliser for $`(q,\mathrm{ad}^{})`$ exists and equals the centre of $`q`$. But $`\widehat{q}=qq^{}`$ is nilpotent and quadratic. Therefore $`(\widehat{q},\mathrm{ad}^{})(\widehat{q},\mathrm{ad})`$ has no generic stabiliser.
$``$ Suppose $`q=g`$ is reductive. Then $`gg^{}`$.
By , $`(g:V^{})`$ always has a generic stabiliser. Assume that this stabiliser is reductive. There is no much harm in it, since there are finitely many $`g`$-modules whose generic stabiliser is not reductive. Then our goal is to prove that the existence of a generic stabiliser for $`(r,\mathrm{ad}^{})`$ imposes a very strong constraint on the action $`(G:V)`$.
Let $`\mathrm{\Omega }_V^{}`$ be the open subset of $`g`$-generic points in $`V^{}`$. Fix a generic stabiliser $`hg`$ and a Cartan subalgebra $`t_hh`$.
5.6 Lemma. There is an open $`R`$-stable subset $`\mathrm{\Xi }(r^{})^{reg}(g\times \mathrm{\Omega }_V^{})`$ such that if $`\eta =(\alpha ,\xi )\mathrm{\Xi }`$, then $`r_\eta `$ is a direct sum of two spaces, one lying in $`g^{}`$ and another lying in $`V^{}`$. Furthermore, eventually replacing $`\eta `$ with an $`H`$-conjugate point, one can achieve that $`r_\eta =t_h(g\xi )^{}`$. Proof. Since $`h`$ is reductive, the $`h`$-modules $`h`$ and $`h^{}`$ can be identified using the restriction to $`h`$ of a non-degenerate $`g`$-invariant symmetric bilinear form on $`g`$. Suppose $`\eta =(\alpha ,\xi )(r^{})^{reg}(g\times \mathrm{\Omega }_V^{})`$. Without loss of generality, assume that $`g_\xi =h`$. As was explained above, the $`r`$-regularity of $`\eta `$ means that $`\alpha `$ is $`h`$-regular as an element of $`h^{}`$. Having identified $`h^{}`$ and $`h`$, we may assume that $`\alpha `$ is regular semisimple. This last condition distinguishes the required subset $`\mathrm{\Xi }`$. Then $`\mathrm{ker}(𝒦_\alpha |_h)`$ is a Cartan subalgebra of $`h`$, and if $`s\mathrm{ker}(𝒦_\alpha |_h)`$, then $`\mathrm{ad}_g^{}(s)\alpha =0`$. Comparing this with Proposition 5, we see that $`r_\eta =\mathrm{ker}(𝒦_\alpha |_h)(g\xi )^{}`$. Taking an $`H`$-conjugate, which does not affect $`\xi `$, we may achieve that $`\mathrm{ker}(𝒦_\alpha |_h)=t_h`$. Thus, for almost all $`r`$-regular points in $`r^{}`$, their stabilisers are conjugate to subalgebras of the form $`\stackrel{~}{h}=t_h(g\xi )^{}`$. Set $`U=(g\xi )^{}`$. By the very construction, $`U`$ is $`h`$-stable. Since $`\eta `$ is regular and therefore $`r_\eta `$ is commutative, $`t_h`$ acts trivially on $`U`$, i.e., $`t_hU=0`$.
5.7 Proposition.
1. Suppose $`(r,\mathrm{ad}^{})`$ has a generic stabiliser. Then $`gUU=\{0\}`$.
2. If $`h=0`$, then the converse is also true. Proof. 1. By Lemma 5 and Eq. (3.1), $`(r,\mathrm{ad}^{})`$ has a generic stabiliser if and only if $`[r,\stackrel{~}{h}]\stackrel{~}{h}=\{0\}`$. We have
$$[r,\stackrel{~}{h}]=[gV,t_hU]=[g,t_hU]+t_hV.$$
Clearly, $`gU`$ is a subspace of $`[g,t_hU]`$. Hence we get the condition that $`gUU=\{0\}`$.
2. Let $`V=UV^{}`$ be an $`h`$-stable decomposition. Then $`t_hV=t_hV^{}V^{}`$. Hence this summand causes no harm. If $`h=0`$, then $`[g,t_hU]=gU`$. Therefore the condition $`gUU=\{0\}`$ appears to be necessary and sufficient for the existence of a generic stabiliser. Recall from the notion of a polar representation of a reductive group. Let $`vV`$ be semisimple, i.e., $`Gv`$ is closed. Define $`c_v=\{xVgxgv\}`$. Then $`(G:V)`$ is said to be polar if there is a semisimple $`vV`$ such that $`dimc_v=dimV//G`$. Such $`c`$ is called a Cartan subspace. Polar representations have a number of nice (and hence restrictive) properties. For instance, all points of $`c`$ are semisimple, all Cartan subspaces are $`G`$-conjugate, the group $`W_c:=𝖭(c)/𝖹(c)`$ is finite, and $`\text{𝕜}[V]^G\text{𝕜}[c]^{W_c}`$ . The latter implies that $`\text{𝕜}[V]^G`$ is polynomial and the morphism $`\pi _G:VV//G`$ is equidimensional .
Our main result related to Question (Q2) is:
5.8 Theorem. Suppose the action $`(G:V)`$ is stable. Then $`(r=gV,\mathrm{ad}^{})`$ has a generic stabiliser if and only if $`(G:V)`$ is a polar representation. Proof. 1. Suppose $`(gV,\mathrm{ad}^{})`$ has a generic stabiliser.
Choose $`\eta =(\alpha ,\xi )\mathrm{\Xi }`$ as prescribed by Lemma 5, so that $`r_\eta =t_hU`$ is a generic stabiliser and hence $`gUU=\{0\}`$ (Proposition 5). In view of stability, we may also assume that $`\xi `$ is ($`g`$-regular and) semisimple. Let us prove that $`U`$ is a Cartan subspace of $`V`$.
As is well known, $`dimV//G=dimV^{}//G`$ and $`(G:V)`$ is stable if and only if $`(G:V^{})`$ is, see e.g. . By the stability hypothesis,
$$\underset{\nu V^{}}{\mathrm{max}}dimG\nu =\underset{vV}{\mathrm{max}}dimGv=dimVdimV//G.$$
Hence $`G\xi =dimVdimV//G`$ and $`dimU=dimV//G`$.
Claim. There is a closed $`G`$-orbit of maximal dimension meeting $`U`$.
Proof of the Claim. The proof of main results in is based on transcendental methods (compact real forms of $`G`$, Kempf–Ness theory). This is an excuse for our using similar methods below. In the next paragraph, we assume that $`\text{𝕜}=`$.
Let $`G_c`$ be a maximal compact subgroup of $`G`$ with Lie algebra $`g_c`$. Fix a $`G_c`$-invariant Hermitian form $`<,>`$ on $`V^{}`$. Without loss of generality, we may assume that $`\xi `$ is of minimal length in $`G\xi `$ and hence $`<g\xi ,\xi >=0`$, see \[11, Sect. 1\]. Upon the identification the $`g_c`$-modules $`V`$ and $`V^{}`$ via $`<,>`$, $`\xi `$ appears to be a point of $`U`$. If $`\stackrel{~}{v}U`$ corresponds to $`\xi `$ under this identification, then we still have $`<g_c\stackrel{~}{v},\stackrel{~}{v}>=0`$, and therefore $`<g\stackrel{~}{v},\stackrel{~}{v}>=0`$. Hence $`G\stackrel{~}{v}`$ is closed \[11, Theorem 1.1\]. Since $`(g_c)_\xi =(g_c)_{\stackrel{~}{v}}`$ and $`(g_c)_{\stackrel{~}{v}}`$ is a compact real form of $`g_{\stackrel{~}{v}}`$ \[11, Prop. 1.3\], we conclude that $`dimg_{\stackrel{~}{v}}=dimg_\xi =dimh`$. $`\mathrm{}`$
For $`\stackrel{~}{v}U`$, we have $`dimg\stackrel{~}{v}=dimVdimU`$. Hence $`g\stackrel{~}{v}=gU`$ for dimension reason. In particular, $`gyg\stackrel{~}{v}`$ for any $`yU`$. Thus, $`U`$ satisfies all conditions in the definition of a Cartan subspace.
2. Suppose $`(G:V)`$ is stable and polar.
Let $`vV`$ be a regular semisimple element and $`c=c_v`$ the corresponding Cartan subspace. Then $`V=gcc`$ and $`gc=gv`$ \[11, Section 2\]. Set $`h=g_v`$. The Lie algebra $`s:=t_hc`$ is commutative, and a direct verification shows that it satisfies Eq. (3.1). Indeed,
$$[gV,t_hc]=[g,t_hc]+t_hV.$$
Using the $`t_h`$-stable decomposition $`V=gcc`$, we see that $`t_hVgc`$. As for the first summand, its $`g`$-component does not belong to $`t_h`$ and its $`V`$-component belongs to $`gc`$. Hence $`[gV,t_hc](t_hc)=\{0\}`$. It remains to find an $`\eta r^{}`$ such that $`r_\eta =s`$.
The dual version of the previous Claim shows that $`(gc)^{}`$ is a Cartan subspace of $`V^{}`$ and that, for sufficiently general $`\xi (gc)^{}`$, we have $`g_\xi =h`$ and $`g\xi =c^{}`$. Now, take an $`\alpha g^{}`$ such that under the identification $`g^{}g`$ it becomes a regular element of $`t_h`$ (i.e., $`\alpha (t_h)^{reg}`$). Then $`\gamma =(\alpha ,\xi )(r^{})^{reg}`$ and $`r_\gamma =s`$. We mention without proof the following consequence of Theorem 5.
5.9 Corollary. If a generic stabiliser $`h`$ for $`(r,\mathrm{ad}^{})`$ is near-toral, then $`\mathrm{rk}h=\mathrm{rk}g`$ and $`U=V^h`$. In case of $`g`$ simple, this implies that $`V`$ is either the adjoint or ”little adjoint” $`g`$-module. (The latter means that the highest weight is the short dominant root, in case $`g`$ has roots of different length.) Remark. It may happen that a generic stabiliser for $`(G:V^{})`$ is not reductive, but $`(gV,\mathrm{ad}^{})`$ still has a generic stabiliser. Indeed, there are $`G`$-modules $`V`$ such that $`r=gV`$ is Frobenius, i.e., $`r^{}`$ has a dense $`R`$-orbit, which certainly ensures the existence of a generic stabiliser. For $`G`$ simple, the list of such $`V`$ is obtained in .
Finally, we consider the field of rational invariants for the coadjoint representation of $`r=gV`$. By ,
$$\mathrm{trdeg}\text{𝕜}(r^{})^R=\mathrm{trdeg}\text{𝕜}(V^{})^G+\mathrm{ind}h,$$
where $`h`$ is a generic stabiliser for $`(G:V^{})`$. It follows from Eq. (5.4) that $`\text{𝕜}(V^{})^G`$ can be regarded as a subfield of $`\text{𝕜}(r^{})^R`$.
5.10 Theorem. If $`\mathrm{ind}h=0`$, then $`\text{𝕜}(r^{})^R\text{𝕜}(V^{})^G`$. Proof. It suffices to verify that $`\text{𝕜}(V^{})^G`$ separates $`R`$-orbits in a dense open subset of $`r^{}`$. Let $`p:r^{}=V^{}g^{}V^{}`$ denote the projection. If $`𝒪V^{}`$ is a generic $`G`$-orbit, then we will prove that $`p^1(𝒪)`$ contains a dense $`R`$-orbit. The latter is equivalent to that, for any $`\xi 𝒪`$, $`G_\xi V`$ has a dense orbit in $`p^1(\xi )=\{\xi \}\times g^{}`$. Since $`1V`$ is a normal subgroup of $`G_\xi V`$, we first look at its orbits. For any $`(\xi ,\alpha )p^1(\xi )`$, we have $`(1V)(\xi ,\alpha )=(\xi ,\alpha +V\xi )`$. Hence all orbits are parallel affine space of dimension $`dim(V\xi )`$. Therefore, it will be sufficient to prove that $`G_\xi `$ has a dense orbit in the (geometric) quotient $`p^1(\xi )/(1V)`$. Because $`V\xi =(g_\xi )^{}`$, that quotient is isomorphic to $`g^{}/(g_\xi )^{}(g_\xi )^{}`$ as $`G_\xi `$-variety. Now, the presence of a dense $`G_\xi `$-orbit in $`(g_\xi )^{}`$ exactly means that $`\mathrm{ind}g_\xi =0`$, which is true as $`\xi `$ is generic. Remarks. 1. In Theorem 5, the reductivity of $`G`$ is not needed. It suffices to assume that $`(G:V^{})`$ has a generic stabiliser.
2. A related result for $`\text{𝕜}(r^{})^R`$ is obtained in \[31, Corollary 2.9\] under the assumption that $`\mathrm{trdeg}\text{𝕜}(V^{})^G=0`$, but without assuming that $`G`$ is reductive.
## 6. Reductive semi-direct products and their polynomial invariants
In this section, we study polynomial invariants of semi-direct products $`q=gV`$, where $`g`$ is reductive.
Our main technical tool is the following result of Igusa (see \[19, Lemma 4\], \[46, Theorem 4.12\]). For reader’s convenience, we provide a proof. Given an irreducible variety $`Y`$, we say that an open subset $`\mathrm{\Omega }Y`$ is big if $`Y\mathrm{\Omega }`$ contains no divisors.
6.1 Lemma (Igusa). Let $`A`$ be an algebraic group acting regularly on an irreducible affine variety $`X`$. Suppose $`S`$ is an integrally closed finitely generated subalgebra of $`\text{𝕜}[X]^A`$ and the morphism $`\pi :X\mathrm{Spec}S=:Y`$ has the properties:
(i) the fibres of $`\pi `$ over a dense open subset of $`Y`$ contain a dense $`A`$-orbit;
(ii) $`\mathrm{Im}\pi `$ contains a big open subset of $`Y`$.
Then $`S=\text{𝕜}[X]^A`$. In particular, the algebra of $`A`$-invariants is finitely generated. Proof. From (i) and Rosenlicht’s theorem, it follows that $`\text{𝕜}(Y)=\text{𝕜}(X)^A`$. In particular, $`\text{𝕜}(X)^A`$ is the quotient field of $`\text{𝕜}[X]^A`$. Assume that $`S\text{𝕜}[X]^A`$. Then one can find a finitely generated intermediate subalgebra: $`S\stackrel{~}{S}\text{𝕜}[X]^A`$ such that $`S\stackrel{~}{S}`$. The natural morphism $`\stackrel{~}{\pi }:\mathrm{Spec}\stackrel{~}{S}Y`$ is birational and its image contains a big open subset of $`Y`$ (because $`\pi `$ does). Since $`Y`$ is normal, the Richardson lemma \[5, 3.2 Lemme 1\] implies that $`\stackrel{~}{\pi }`$ is an isomorphism. This contradiction shows that $`S=\text{𝕜}[X]^A`$. Recall that $`Q:=GV`$ is a connected group with Lie algebra $`q`$. Here $`1V`$ is exactly the unipotent radical of $`Q`$, which is also denoted $`Q^u`$. Let $`T`$ be a maximal torus of $`G`$ with the corresponding Cartan subalgebra $`t`$.
First, we consider the adjoint representation of $`gV`$.
6.2 Theorem. Let $`V`$ be an arbitrary $`G`$-module, $`q=gV`$, and $`m=dimV^t`$. Then
* $`\text{𝕜}[q]^{Q^u}`$ is a polynomial algebra of Krull dimension $`dimg+m`$. It is freely generated by the coordinates on $`g`$ and the functions $`\widehat{F}_i`$, $`i=1,\mathrm{},m`$, associated with covariants of type $`V^{}`$ (see below).
* $`\text{𝕜}[q]^Q`$ is a polynomial algebra of Krull dimension $`dimt+m`$. It is freely generated by the basic invariants of $`\text{𝕜}[g]^G`$ and the same functions $`\widehat{F}_i`$, $`i=1,\mathrm{},m`$.
* $`\mathrm{max}dim_{xq}Qx=dimqdimq//Q`$;
* If $`\pi :qq//Q`$ is the quotient morphism, then $`\mathrm{\Omega }:=\{xq\text{d}\pi _x\text{ is onto }\}`$ is a big open subset of $`q`$.
Proof. (i) By Theorem 4, $`\mathrm{Mor}_G(g,V^{})`$ is a free $`\text{𝕜}[g]^G`$-module of rank $`m`$. Let $`F_1,\mathrm{},F_m`$ be a basis for this module and $`\widehat{F}_1,\mathrm{},\widehat{F}_m`$ the corresponding $`Q`$-invariants on $`q`$, i.e., $`\widehat{F}_i(x+ϵv)=F_i(x),v`$. To prove that $`\text{𝕜}[q]^{Q^u}`$ is freely generated by the coordinate functions on $`g`$ and the polynomials $`\widehat{F}_i`$, $`i=1,\mathrm{},m`$, we wish to apply Lemma 6.
Set $`X_m=\{xgdim\text{span}\{F_1(x),\mathrm{},F_m(x)\}=m\}`$. That is, $`X_m`$ is the set of those $`x`$, where the vectors $`F_i(x)V^{}`$, $`i=1,\mathrm{},m`$, are linearly independent.
Claim. $`X_m`$ is a big open subset of $`g`$. More precisely, $`\mathrm{codim}_g(gX_m)3`$.
Proof of the claim. The set of regular elements of $`g`$, $`g^{reg}`$, has the property that $`\mathrm{codim}(gg^{reg})3`$ and $`\overline{Gx}`$ is normal for any $`xg^{reg}`$ . The condition that $`\mathrm{codim}_{\overline{Gx}}(\overline{Gx}Gx)2`$ is satisfied for every $`xg`$, since any $`G`$-orbit is even-dimensional. By Theorems 4 and 4, we conclude that $`X_mg^{reg}`$, and the claim follows.
Let $`x_1,\mathrm{},x_n`$ be the coordinates on $`g`$, where $`n=dimg`$. Then $`x_1,\mathrm{},x_n,\widehat{F}_1,\mathrm{},\widehat{F}_m`$ are algebraically independent, because their differentials are linearly independent on $`X_mV`$. Consider the mapping
$$\tau :q\mathrm{Spec}\text{𝕜}[x_1,\mathrm{},x_n,\widehat{F}_1,\mathrm{},\widehat{F}_m]=\text{𝕜}^{n+m},$$
where $`\tau (x+ϵv)=(x,\widehat{F}_1(x+ϵv),\mathrm{},\widehat{F}_m(x+ϵv))`$. We identify $`\text{𝕜}^{n+m}`$ with $`g\times \text{𝕜}^m`$. If $`x=(x_1,\mathrm{},x_n)X_m`$, then the $`F_i(x)`$’s are linearly independent, so that the system
$$\widehat{F}_i(x+ϵv)=F_i(x),v=\alpha _i,i=1,\mathrm{},m$$
has a solution $`v`$ for any $`m`$-tuple $`\alpha =(\alpha _1,\mathrm{},\alpha _m)`$. Hence $`\mathrm{Im}\tau X_m\times \text{𝕜}^m`$, which means that $`\mathrm{Im}\tau `$ contains a big open subset of $`\text{𝕜}^{n+m}`$.
It follows from the above Claim that $`g^{reg}=X_m\{ygdimGy=nm\}`$. Take $`xg^{reg}`$, and let $`v_\alpha `$ be a solution to the system $`\widehat{F}_i(x+ϵv)=\alpha _i`$. Then $`\tau ^1(x,\alpha )x+ϵv_\alpha `$ and
$$\tau ^1(x,\alpha )Q^u(x+ϵv_\alpha )=\{x+ϵ(v_\alpha +xV)\}.$$
Since $`xX_m`$, we have $`dim\tau ^1(x,\alpha )=nm`$. On the other hand, $`dim[g,x]=nm`$, by the definition of $`g^{reg}`$. Hence $`\tau ^1(x,\alpha )=Q^u(x+ϵv_\alpha )`$ for dimension reason. Thus, a generic fibre of $`\tau `$ is a $`Q^u`$-orbit, and Lemma 6 applies here.
(ii) Clearly,
$$\text{𝕜}[q]^Q=(\text{𝕜}[q]^{Q^u})^G=\text{𝕜}[x_1,\mathrm{},x_n,\widehat{F}_1,\mathrm{},\widehat{F}_m]^G.$$
Since the $`\widehat{F}_i`$’s are already $`G`$-invariant, the algebra in question is equal to
$$\text{𝕜}[g]^G[\widehat{F}_1,\mathrm{},\widehat{F}_m].$$
(iii) The dimension of a $`Q`$-orbit cannot be greater than $`dimqdimq//Q`$, and if $`xt`$ is regular, then $`dimQ(x+ϵ0)=dimQdimtm`$.
(iv) It follows from the previous discussion that $`\mathrm{\Omega }g^{reg}V`$. 6.3 Remarks. 1. If $`V^T=\{0\}`$, then the module of covariants of type $`V^{}`$ is trivial, so that we obtain a natural isomorphism $`\text{𝕜}[q]^Q\text{𝕜}[g]^G`$.
2. From Theorem 5 and Proposition 5, it follows that $`\stackrel{~}{t}:=tV^t`$ is a generic centraliser in $`q`$ and $`\stackrel{~}{W}=N_Q(\stackrel{~}{t})/Z_Q(\stackrel{~}{t})`$ is isomorphic to $`W=N_G(t)/Z_G(t)`$, the usual Weyl group of $`g`$. Therefore
$$\text{𝕜}(q)^Q\text{𝕜}(tV^t)^W=\text{𝕜}(t\times V^t)^W.$$
Since $`\text{𝕜}(t)^W`$ is a rational field, Theorem 5(iv) implies that $`\text{𝕜}(q)^Q`$ is rational, too. For $`g`$ semisimple, the rationality of $`\text{𝕜}(q)^Q`$ also follows from Theorem 6, because in this situation $`\text{𝕜}(q)^Q`$ is the quotient field of $`\text{𝕜}[q]^Q`$. However, if $`V^t0`$, then the restriction homomorphism
$$res:\text{𝕜}[q]^Q\text{𝕜}[t\times V^t]^W$$
is not onto. For, the description of the generators of $`\text{𝕜}[q]^Q`$ shows that $`\text{𝕜}[V^t]^W`$ does not belong to the image of $`res`$. Now, we look at polynomial invariants of the coadjoint representation of $`q=gV`$. As we know from Section 5, the existence of a generic stabiliser for $`(q,\mathrm{ad}^{})`$ is a rare phenomenon; but this existence is not always needed for describing invariants. It follows from Eq. (5.4) that $`\text{𝕜}[V^{}]^G`$ can be regarded as a subalgebra of $`\text{𝕜}[q^{}]^Q`$. Recall that $`\mathrm{trdeg}\text{𝕜}(q^{})^Q=\mathrm{trdeg}\text{𝕜}(V^{})^G+\mathrm{ind}h`$, where $`h`$ is a generic stabiliser for $`(G:V^{})`$. In particular, if $`g`$ is semisimple and $`h`$ is reductive, then $`\mathrm{trdeg}\text{𝕜}(q^{})^Q=\mathrm{trdeg}\text{𝕜}(V^{})^G+\mathrm{rk}h`$. Since the roles of $`V`$ and $`g`$ are interchanged in the dual space, one might hope that $`\text{𝕜}[q^{}]^Q`$ could be generated by $`\text{𝕜}[V^{}]^G`$ and certain invariants arising from $`\mathrm{Mor}_G(V,g^{})`$. This is however false, because it can happen that $`\mathrm{rk}h>0`$, but $`\mathrm{Mor}_G(V,g^{})=0`$. In general, it is not clear how to discover ”missing” invariants associated with the summand $`\mathrm{ind}h`$ (or $`\mathrm{rk}h`$).
The simplest case is that in which $`h=0`$. Then we are in a position to state an analogue of Theorem 6.
6.4 Theorem. As above, let $`q=g\times V`$ and $`Q^u=1V`$. Suppose a generic stabiliser for $`(G:V^{})`$ is trivial. Then $`\text{𝕜}[q^{}]^{Q^u}=\text{𝕜}[V^{}]`$ and $`\text{𝕜}[q^{}]^Q=\text{𝕜}[V^{}]^G`$. Proof. The second equality stems from the first. To prove the first equality, we use the same method as in Theorem 6. The natural projection $`q^{}q^{}/g^{}V^{}`$ is $`Q^u`$-equivariant and satisfies all the requirements of Lemma 6. The details are left to the reader. Remark. In Theorem 6, the reductivity of $`G`$ is not needed.
## 7. Takiff Lie algebras and their invariants
For $`g`$ semisimple, some interesting results on the invariants of $`(gg,\mathrm{ad})`$ are obtained by Takiff in . For this reason, Lie algebras of the form $`qq`$ are sometimes called Takiff (Lie) algebras, see ,. We will follow this terminology.
In this section, we consider orbits and invariants of certain representations of a Takiff group $`\widehat{Q}=Qq`$. Some results on rational invariants have already appeared in Section 5. Our main object here is the polynomial (regular) invariants. We obtain a generalisation of the main result in , which concerns several aspects. First, in place of semisimple Lie algebras, we consider a wider class. Second, the initial representation of $`Q`$ is not necessarily adjoint. Third, we also describe the invariants of the unipotent group $`1q\widehat{Q}`$. Fourth, our proof does not exploit complex numbers and complex topology.
If $`V`$ is a $`q`$-module, then $`V\times V`$ can regarded as $`qq`$-module in a very natural way. For $`(x_1,x_2)qq`$ and $`(v_1,v_2)V\times V`$, we define
$$(x_1,x_2)(v_1,v_2):=(x_1v_1,x_1v_2x_2v_1).$$
This $`q`$-module will be denoted by $`\widehat{V}=VV`$. We also write $`v_1+ϵv_2`$ for $`(v_1,v_2)`$. If $`f\text{𝕜}[V]^Q`$, then $`\text{d}f\mathrm{Mor}_Q(V,V^{})`$, and we define $`\widehat{F}_f\text{𝕜}[\widehat{V}]`$ by the rule: $`\widehat{F}_f(x+ϵy)=\text{d}f_x,y`$. Similarly to Lemma 4, one proves
$$\widehat{F}_f\text{𝕜}[\widehat{V}]^{\widehat{Q}}.$$
Here one needs the fact that $`\text{d}f_v`$ annihilates the tangent space of $`Qv`$ at $`vV`$.
7.1 Theorem. Let $`V`$ be a $`Q`$-module. Suppose the action $`(Q:V)`$ satisfies the following conditions:
(1) $`\text{𝕜}[V]^Q`$ is a polynomial algebra;
(2) $`\mathrm{max}dim_{vV}Qv=dimVdimV//Q`$;
(3) If $`\pi _Q:VV//Q`$ is the quotient morphism and $`\mathrm{\Omega }:=\{vV(\text{d}\pi _Q)_v\text{ is onto }\}`$, then $`V\mathrm{\Omega }`$ contains no divisors.
Then
* $`\text{𝕜}[\widehat{V}]^{1q}`$ is a polynomial algebra of Krull dimension $`dimV+dimV//Q`$, which is generated by the coordinates on the first factor of $`\widehat{V}`$ and the polynomials $`\widehat{F}_1,\mathrm{},\widehat{F}_m`$ associated with the differentials of basic invariants in $`\text{𝕜}[V]^Q`$;
* $`\text{𝕜}[\widehat{V}]^{\widehat{Q}}`$ is a polynomial algebra of Krull dimension 2$`dimV//Q`$, which is freely generated by the basic invariants of $`\text{𝕜}[V]^Q`$ and the same functions $`\widehat{F}_i`$, $`i=1,\mathrm{},m`$.
* The $`\widehat{Q}`$-module $`\widehat{V}`$ satisfies conditions (1)–(3), too.
Proof. The proof is very close in the spirit to the proof of Theorem 6, though some technical details are different.
Set $`N=1q`$. Let $`f_1,\mathrm{},f_m`$, $`m=dimV//Q`$, be algebraically independent generators of $`\text{𝕜}[V]^Q`$. As was noticed above, to each $`f_i`$ one may associate the polynomial $`\widehat{F}_i=\widehat{F}_{f_i}\text{𝕜}[VV]^{\widehat{Q}}`$.
(i) We are going to prove, using Lemma 6, that $`\text{𝕜}[\widehat{q}]^N`$ is freely generated by the coordinate functions on $`V`$ (which is the first component of $`\widehat{V}`$) and the polynomials $`\widehat{F}_i`$, $`i=1,\mathrm{},m`$. Let $`x_1,\mathrm{},x_n`$ be the coordinate functions on $`V`$. Then $`x_1,\mathrm{},x_n,\widehat{F}_1,\mathrm{},\widehat{F}_m`$ are algebraically independent, because their differentials are linearly independent on $`\mathrm{\Omega }V`$. Consider the mapping
$$\widehat{\tau }:\widehat{V}\mathrm{Spec}\text{𝕜}[x_1,\mathrm{},x_n,\widehat{F}_1,\mathrm{},\widehat{F}_m]=\text{𝕜}^{n+m}.$$
We identify $`\text{𝕜}^{n+m}`$ with $`V\times \text{𝕜}^m`$. If $`x=(x_1,\mathrm{},x_n)\mathrm{\Omega }`$, then $`(\text{d}f_i)_x`$ are linearly independent, so that the system $`\widehat{F}_i(x+ϵy)=\alpha _i`$, $`i=1,\mathrm{},m`$, has a solution $`y`$ for any $`m`$-tuple $`\alpha =(\alpha _1,\mathrm{},\alpha _m)`$. Hence $`\mathrm{Im}\widehat{\tau }\mathrm{\Omega }\times \text{𝕜}^m`$, which means that $`\mathrm{Im}\widehat{\tau }`$ contains a big open subset of $`\text{𝕜}^{n+m}`$.
Next, consider $`\mathrm{\Omega }^{}=\mathrm{\Omega }\{yVdimQy=nm\}`$. In view of condition (2), it is still a non-empty open $`Q`$-stable subset of $`V`$. Take $`x\mathrm{\Omega }^{}`$, and let $`y_\alpha `$ be a solution to the system $`\widehat{F}_i(x+ϵy)=(\text{d}f_i)_x,y=\alpha _i`$, $`i=1,\mathrm{},m`$. Then $`\widehat{\tau }^1(x,\alpha )x+ϵy_\alpha `$ and
$$\widehat{\tau }^1(x,\alpha )N(x+ϵy_\alpha )=\{x+ϵ(y_\alpha +qx)\}.$$
Since $`x\mathrm{\Omega }`$, we have $`dim\widehat{\tau }^1(x,\alpha )=nm`$. On the other hand, $`dimQx=nm`$, because of the definition of $`\mathrm{\Omega }^{}`$. Hence $`\widehat{\tau }^1(x,\alpha )=N(x+ϵy_\alpha )`$ for dimension reason. Thus, a generic fibre of $`\widehat{\tau }`$ is an $`N`$-orbit, and Lemma 6 applies here.
(ii) Clearly,
$$\text{𝕜}[\widehat{V}]^{\widehat{Q}}=(\text{𝕜}[\widehat{V}]^N)^Q=\text{𝕜}[x_1,\mathrm{},x_n,\widehat{F}_1,\mathrm{},\widehat{F}_m]^Q.$$
Since the $`\widehat{F}_i`$’s are already $`Q`$-invariant, the algebra in question is equal to
$$\text{𝕜}[V]^Q[\widehat{F}_1,\mathrm{},\widehat{F}_m]=\text{𝕜}[f_1,\mathrm{},f_m,\widehat{F}_1,\mathrm{},\widehat{F}_m].$$
(iii) We have to check that the $`\widehat{Q}`$-module $`\widehat{V}`$ satisfies properties (1)–(3).
$``$ Property (1) is verified in (ii).
$``$ If $`x\mathrm{\Omega }`$, then $`dim\widehat{Q}(x+ϵ0)=2n2m`$, which gives property (2) for $`\widehat{Q}`$.
$``$ Set $`\widehat{\mathrm{\Omega }}=\mathrm{\Omega }\times V`$. It is a big open subset of $`\widehat{V}`$. Explicit expressions for algebraically independent generators of $`\text{𝕜}[\widehat{V}]^{\widehat{Q}}`$ show that their differentials are linearly independent on $`\widehat{\mathrm{\Omega }}`$, which is exactly Property (3) for $`\widehat{V}`$. 7.2 Remarks. 1. If the pair $`(q,V)`$ satisfies properties (1)-(3) of Theorem 7, then the passage $`(q,V)(\widehat{q},\widehat{V})`$ can be iterated ad infinitum without losing those properties.
2. Since the adjoint representation of a semisimple Lie algebra $`g`$ has properties (1)–(3), iterating the Takiffisation procedure $`ggg`$ always yields algebras with a polynomial ring of invariants for the adjoint representation. This is the main result of . Explicit form of the basic invariants for $`(gg,\mathrm{ad})`$ is also pointed out there. Notice also that Takiff’s results follow from either Theorem 6 with $`V=g`$ or Theorem 7 with $`q=V=g`$.
3. By Theorem 6, the adjoint representation of $`q=gV`$ satisfies all the conditions of Theorem 7. Therefore these $`q`$ can be used as building blocks for Takiffisation procedure, which yields more and more complicated Lie algebras having polynomial algebras of invariants. Let us make some comments on the conditions of Theorem 7. If $`Q=G`$ is semisimple, then conditions (2) and (3) are always satisfied, regardless of the fact whether $`\text{𝕜}[V]^G`$ is polynomial. For condition (3) we refer to \[20, Satz 2\], while (2) follows since $`G`$ has no rational characters and therefore $`\text{𝕜}(V)^G`$ is the quotient field of $`\text{𝕜}[V]^G`$. Thus, we have
7.3 Corollary. If $`\rho :GGL(V)`$ is a representation of a semisimple group such that $`\text{𝕜}[V]^G`$ is polynomial, then Theorem 7 applies to the $`Gg`$-module $`VV`$.
## 8. The null-cone and its irreducibility
In previous sections, we described several instances of representations of nonreductive Lie algebras having a polynomial algebra of invariants. If $`QGL(\stackrel{~}{V})`$ and $`\text{𝕜}[\stackrel{~}{V}]^Q`$ is polynomial, then it is natural to inquire of whether it is true that $`\text{𝕜}[\stackrel{~}{V}]`$ is a free $`\text{𝕜}[\stackrel{~}{V}]^Q`$-module. As is well known, the freeness is equivalent to that the quotient morphism $`\pi :\stackrel{~}{V}\stackrel{~}{V}//Q`$ is equidimensional, i.e., has the property that $`dim\pi ^1(\pi (0))=dim\stackrel{~}{V}dim\stackrel{~}{V}//Q`$. As in the case of reductive group actions, we say that $`\pi ^1(\pi (0))`$ is the null-cone, denoted $`N^Q(\stackrel{~}{V})`$ or $`N(\stackrel{~}{V})`$.
In this section, we only deal with reductive semi-direct products and their representations. Our goal is to describe necessary and sufficient conditions for equidimensionality of $`\pi `$ and point out some consequences of it. We consider two types of representations:
A) $`q=gV`$, where $`V`$ is a $`g`$-module, and $`\stackrel{~}{V}=q`$, i.e., we consider the adjoint representation of $`q`$.
B) $`q=gg`$ is a reductive Takiff algebra and $`\stackrel{~}{V}=VV`$, where $`V`$ is a $`g`$-module.
We begin with case A). Recall that $`m=dimV^T`$ and $`F_1,\mathrm{},F_m`$ is a basis for the $`\text{𝕜}[g]^G`$-module $`\mathrm{Mor}_G(g,V^{})`$. The null-cone for $`(g,\mathrm{ad})`$ is denoted by $`𝒩(g)`$ or merely by $`𝒩`$. In other words, $`𝒩`$ is the set of nilpotent elements of $`g`$. Recall that $`𝒩`$ is irreducible and $`dim𝒩=dimgdimg//G=dimgdimt`$.
Theorem 6 says that if $`V^T=0`$, then $`\text{𝕜}[q]^Q=\text{𝕜}[g]^G`$ and therefore $`N(q)𝒩\times V`$. In this trivial case, $`\pi _Q`$ is equidimensional, since it is so for $`\pi _G:gg//G`$. Therefore we assume below that $`V^T0`$.
Define a stratification of $`g`$ in the following way:
$$X_{i,V}=X_i=\{xgdim\text{span}\{F_1(x),\mathrm{},F_m(x)\}=i\}.$$
Then $`\overline{X}_i\overline{X}_{i+1}`$ and $`\overline{X}_m=g`$. The induced stratification on the null-cone is $`X_i(𝒩):=X_i𝒩`$. As is shown in the proof of Theorem 6, $`X_m`$ is a big open subset of $`g`$ containing $`g^{reg}`$. Therefore $`X_m(𝒩)`$ is a big open subset of $`𝒩`$ containing the regular nilpotent orbit.
8.1 Theorem.
* The quotient morphism $`\pi _Q:qq//Q`$ is equidimensional if and only if
$`\mathrm{codim}_𝒩X_i(𝒩)mi`$.
* If $`N(q)`$ is irreducible, then $`\pi _Q`$ is equidimensional;
* $`N(q)`$ is irreducible if and only if $`\mathrm{codim}_𝒩X_i(𝒩)mi+1`$ for $`i<m`$.
Proof. 1. Since $`\pi _Q`$ is dominant, all irreducible components of $`N(q)`$ are of dimension $`dimqdimq//Q`$. By Theorem 6, $`N(q)=\{(x,v)x𝒩\&\widehat{F}_i(x,v)=F_i(x),v=0i\}`$. Let $`p:N(q)𝒩`$ be the projection onto the first factor. Then $`N(q)={\displaystyle \underset{i=0}{\overset{m}{}}}p^1(X_i(𝒩))`$ and $`dimp^1(X_i(𝒩))=dimX_i(𝒩)+dimVi`$.
2. By Theorem 6, if $`e𝒩^{reg}`$, then $`(e,0)q^{reg}`$ and $`(\text{d}\pi _Q)_{(e,0)}`$ is onto . Therefore, $`(e,0)`$ is a smooth point of $`N(q)`$, and the unique irreducible component of $`N(q)`$ to which $`(e,0)`$ belongs is of dimension $`dimqdimq//Q`$. On the other hand, $`dimp^1(X_m(𝒩))=dim𝒩+dimVm=dimqdimq//Q`$. Hence $`\overline{p^1(X_m(𝒩))}`$ is the irreducible component of $`N(q)`$ containing $`(e,0)`$.
3. The proof of part 2 shows that $`\overline{p^1(X_m(𝒩))}`$ is an irreducible component of $`N(q)`$ of expected dimension. To ensure the irreducibility, we have to require that $`\overline{p^1(X_i(𝒩))}`$ cannot be an irreducible component for $`i<m`$. Since all irreducible components of $`N(q)`$ are of dimension $`dimqdimq//Q`$, the condition that $`dimp^1(X_i(𝒩))<dimqdimq//Q`$ for $`i<m`$ is equivalent to the irreducibility. The following is now immediate.
8.2 Corollary. If $`m=1`$, then $`N(q)`$ is irreducible; if $`m=2`$, then $`\pi _Q`$ is equidimensional. 8.3 Remarks. 1. Since $`𝒩`$ consists of finitely many $`G`$-orbits, condition 8(1) is equivalent to the following: if $`GxX_i(𝒩)`$, then $`dimGxdim𝒩(mi)`$, or
$$dimz_g(x)\mathrm{rk}gmdim\left(\text{span}\{F_1(x),\mathrm{},F_m(x)\}\right).$$
Furthermore, a more careful look at the projection $`N(q)𝒩`$ shows that if last condition is satisfied, then the number of the irreducible components of $`N(q)`$ equals the number of the $`G`$-orbits $`Gx𝒩`$ such that $`dimz_g(x)\mathrm{rk}g=mdim(\text{span}\{F_1(x),\mathrm{},F_m(x)\})`$.
2. The condition in Theorem 8(1) for $`i=0`$ reads $`dim𝒩dim(X_0(𝒩))m`$, or $`dimV^Tdim𝒩dimX_0(𝒩)dim𝒩`$. This is a rough necessary condition for $`\pi _Q`$ to be equidimensional. Let $`G`$ be simple and $`V_\lambda `$ a simple $`G`$-module with highest weight $`\lambda `$. Then $`(V_\lambda )^T0`$ if and only if $`\lambda `$ lies in the root lattice, $``$. The function $`ndim(V_{n\lambda })^T`$, $`\lambda `$, has a polynomial growth. The only case in which this function is constant is that of $`G=SL_p`$, $`\lambda =p\phi _1`$ or $`p\phi _{p1}`$. Here $`\phi _i`$’s are fundamental weights, and $`dim(V_{n\lambda })^T=1`$ for any $`n`$. Thus, modulo this exception, there are finitely many simple $`G`$-modules $`V`$ such that $`V^T0`$ and $`\pi _Q`$ is equidimensional. For future use, we record a relationship between the stratifications of $`𝒩`$ and $`g`$.
8.4 Proposition. If $`\mathrm{codim}_𝒩X_i(𝒩)mi+1`$ for $`i<m`$, then $`\mathrm{codim}_gX_imi+2`$. Proof. It follows from the definitions that $`\overline{X_i}\underset{ji}{}X_j`$ and $`\overline{X_i(𝒩)}\underset{ji}{}X_j(𝒩)`$. For $`i<m`$, we have $`X_ig^{reg}=\text{}`$ and hence $`\pi _G(\overline{X_i})=\overline{X_i}//G`$ is a proper subvariety of $`g//G`$. Therefore
$`dim\overline{X_i}dimg//G1+dim\overline{X_i(𝒩)}dimg(mi+2)`$. There is another interesting cone related to $`q=gV`$. Consider the morphism $`\overline{\pi }:q\text{𝕜}^m`$, $`(x,v)(\widehat{F}_1(x,v),\mathrm{},\widehat{F}_m(x,v))`$. The zero-fibre of $`\overline{\pi }`$ is denoted by $`N^u(q)`$. Thus,
$$N^u(q)=\{(x,v)qF_i(x),v=0i=1,\mathrm{},m\}.$$
The proof of the following result is entirely similar to that of Theorem 8. One should only consider the projection $`N^u(q)g`$.
8.5 Theorem.
* The morphism $`\overline{\pi }:q\text{𝕜}^m`$ is equidimensional if and only if $`\mathrm{codim}_gX_imi`$.
* If $`N^u(q)`$ is irreducible, then $`\overline{\pi }`$ is equidimensional;
* $`N^u(q)`$ is irreducible if and only if $`\mathrm{codim}_gX_imi+1`$ for $`i<m`$.
Now, comparing Theorem 8(iii), Proposition 8, and Theorem 8(iii), one concludes that if $`N(q)`$ is irreducible, then so is $`N^u(q)`$.
But one can derive a much stronger assertion on $`N^u(q)`$ from the irreducibility of $`N(q)`$. This is related to properties of symmetric algebras of certain modules over polynomial rings and exploits some technique from , .
Let $`\mathrm{Mor}(g,V^{})`$ be the $`\text{𝕜}[g]`$-module of all polynomial morphisms $`F:gV^{}`$. Consider the homomorphism $`\widehat{\tau }:\mathrm{Mor}(g,V^{})\mathrm{Mor}(g,V^{})`$ defined by $`\widehat{\tau }(F)(x)=xF(x)`$. (Here “$``$” refers to the $`g`$-module structure on $`V^{}`$.)
8.6 Theorem. $`\mathrm{ker}\widehat{\tau }`$ is a free $`\text{𝕜}[g]`$-module of rank $`m`$. More precisely, $`(F_1,\mathrm{},F_m)`$ is a basis for $`\mathrm{ker}\widehat{\tau }`$. Proof. The proof is based on the same idea as the proof of Theorem 1.9 in .
Clearly, $`\mathrm{ker}\widehat{\tau }`$ is a torsion-free $`\text{𝕜}[g]`$-module and the rank $`\mathrm{rk}(\mathrm{ker}\widehat{\tau }):=dim(\mathrm{ker}\widehat{\tau }_{\text{𝕜}[g]}\text{𝕜}(g))`$ is well-defined. An easy argument shows that the rank of $`\widehat{\tau }`$ over $`\text{𝕜}(g)`$ equals $`dimV\mathrm{max}_{xg}dim(V^{})^x=dimVm`$. Hence $`\mathrm{rk}(\mathrm{ker}\widehat{\tau })=m`$. Obviously, $`F_i\mathrm{ker}\widehat{\tau }`$ and $`{\displaystyle \underset{i=1}{\overset{m}{}}}\text{𝕜}[g]F_i`$ is a free submodule of $`\mathrm{ker}\widehat{\tau }`$ of rank $`m`$. It follows that, for any $`F\mathrm{ker}\widehat{\tau }`$, there exist $`\widehat{p},p_1,\mathrm{},p_m\text{𝕜}[g]`$ such that
$$\widehat{p}F=\underset{i}{}p_iF_i.$$
Assume $`\widehat{p}\text{𝕜}^{}`$. Let $`p`$ be a prime factor of $`\widehat{p}`$ and $`D`$ the divisor of zeros of $`p`$. Then $`_ip_i(v)F_i(v)=0`$ for any $`vD`$. Since $`g^{reg}`$ is big, $`g^{reg}D`$ is dense in $`D`$. Because $`\{F_i(v)\}`$ are linearly independent for $`vg^{reg}`$, we obtain $`p_i|_D0`$. Hence $`p_i/p\text{𝕜}[g]`$ for each $`i`$, and we are done. Let $`E`$ denote the $`\text{𝕜}[g]`$-module $`\mathrm{Im}\widehat{\tau }`$. In view of the previous theorem, we have the exact sequence
(8.7)
$$0\underset{i=1}{\overset{m}{}}\text{𝕜}[g]F_i\stackrel{\widehat{\beta }}{}\mathrm{Mor}(g,V^{})\stackrel{\widehat{\tau }}{}E0.$$
Choose a basis $`\xi _1,\mathrm{},\xi _n`$ for $`V^{}`$. Using this basis, we identify $`\mathrm{Mor}(g,V^{})=\text{𝕜}[g]V^{}`$ with $`\text{𝕜}[g]^n`$. Then we can write $`F_j(x)=_{i=1}^nF_{ij}(x)\xi _i`$, where $`F_{ij}\text{𝕜}[g]`$. If we regard sequence (8.7) as a sequence
$$0\text{𝕜}[g]^m\stackrel{\widehat{\beta }}{}\text{𝕜}[g]^n\stackrel{\widehat{\tau }}{}E0,$$
then $`\widehat{\beta }`$ becomes an $`n\times m`$-matrix with entries $`F_{ij}`$. Let $`I_t(\widehat{\beta })`$ be the ideal generated by $`t\times t`$ minors of $`\widehat{\beta }`$. For $`d`$, consider the following condition
$`(_d)`$ $`\mathrm{ht}I_t(\widehat{\beta })mt+1+d`$ for $`1tm`$.
The ideals $`I_t(\widehat{\beta })`$ are independent of the presentation of $`E`$. These are Fitting ideals of $`E`$, see e.g. \[43, 1.1\]. Let $`\text{Sym}_{\text{𝕜}[g]}(E)`$ denote the symmetric algebra of the $`\text{𝕜}[g]`$-module $`E`$.
8.8 Theorem. Suppose $`N(q)`$ is irreducible. Then
* The condition $`(_2)`$ is satisfied by $`E`$.
* $`\text{Sym}_{\text{𝕜}[g]}(E)`$ is a factorial domain of Krull dimension $`dimg+nm`$.
* $`N^u(q)`$ is an irreducible factorial complete intersection, and $`\text{𝕜}[N^u(q)]=\text{Sym}_{\text{𝕜}[g]}E`$.
* $`N^u(q)=\overline{\mathrm{Im}(\kappa )}`$, where $`\kappa :qq`$ is defined by $`\kappa (x,v)=(x,xv)`$, $`xg,vV`$.
Proof. (i) It is easily seen that $`X_i`$ is the zero locus of $`I_{i+1}(\widehat{\beta })`$. Therefore condition $`(_2)`$ is satisfied in view of Proposition 8.
(ii) The exact sequence (8.7) shows that $`E`$ has projective dimension at most one. Therefore part (ii) follows from (i) combined with \[2, Prop. 3 & 6\].
(iii) The universal property of symmetric algebras implies that $`\text{Sym}_{\text{𝕜}[g]}(E)`$ is the quotient of $`\text{Sym}_{\text{𝕜}[g]}(\text{𝕜}[g]V^{})=\text{𝕜}[g\times V]`$ by the ideal generated by the image of $`\widehat{\beta }`$. It follows from the construction that $`\widehat{\beta }(F_i)=\widehat{F}_i`$. Hence $`\text{Sym}_{\text{𝕜}[g]}(E)=\text{𝕜}[N^u(q)]`$, and the other assertions follow from (ii).
(iv) Clearly, $`\overline{\mathrm{Im}(\kappa )}`$ is an irreducible subvariety of $`q`$. Taking the (surjective) projection to $`g`$ and looking at the dimension of the generic fibre, one finds that $`dim\overline{\mathrm{Im}(\kappa )}=dimg+nm`$. Thus, $`\overline{\mathrm{Im}(\kappa )}N^u(q)`$, both have the same dimension and are irreducible. Hence they are equal. Remark. For $`V=g`$, i.e., for the Takiff algebra $`gg`$, condition $`(_2)`$ can be proved directly, without referring to the irreducibility of $`N(q)`$, see \[22, Prop. 2.1\]. In this special case, the above results for $`N^u(q)`$ are already obtained in \[22, Prop. 2.4\]. Actually, $`N(q)`$ is irreducible if $`V=g`$. But this fact, as well as “Takiff” terminology, was not used in loc. cit. In Section 9, we give new examples of semi-direct products $`q=gV`$ such that $`N(q)`$ is irreducible and thereby new instances, where Theorem 8 applies.
Now, we proceed to case B).
Recall that $`\widehat{G}=Gg`$ and $`\widehat{V}=VV`$ is a $`\widehat{G}`$-module. To a great extent, our results in this case are similar to those in case A). A notable distinction is, however, that whereas the adjoint representation of $`G`$ has some good properties for granted, we have to require these properties for $`(G:V)`$.
We will assume below that $`(G:V)`$ satisfies properties (1)–(3) of Theorem 7, with $`G`$ in place of $`Q`$, and use the respective notation. In particular, $`\widehat{G}^u=1V`$, $`m=dimV//G`$, $`\text{𝕜}[V]^G=\text{𝕜}[f_1,\mathrm{},f_m]`$, and $`\widehat{F}_i\text{𝕜}[\widehat{V}]^{\widehat{G}}`$ is the invariant associated with $`\text{d}f_i`$. As in case A), we define a stratification of $`V`$ by
$$Y_i=\{xVdim\text{span}\{(\text{d}f_1)_x,\mathrm{},(\text{d}f_m)_x\}=i\}=\{xV\mathrm{rk}(\text{d}\pi _G)_x=i\}.$$
Then $`\overline{Y}_i\overline{Y}_{i+1}`$ and $`\overline{Y}_m=V`$. Notice that $`Y_0=\{0\}`$. The induced stratification of the null-cone $`N^G(V)=N(V)`$ is $`Y_i(N(V)):=Y_iN(V)`$. Since $`\pi _G:VV//G`$ is onto, $`dimN(V)dimVdimV//G`$. But, unlike the case of $`(G:g)`$, it may happen that the last inequality is strict and $`Y_m(N(V))=\text{}`$.
8.9 Lemma. Suppose $`\pi _{\widehat{G}}:\widehat{V}\widehat{V}//\widehat{G}`$ is equidimensional. Then so is $`\pi _G:VV//G`$ and $`Y_m(N(V))\text{}`$. Proof. Consider the projection $`p:N(\widehat{V})N(V)`$. If $`j`$ is the maximal index such that $`Y_j(N(V))\text{}`$, then $`dimN(\widehat{V})=dimN(V)+dimVj`$. Since $`dimN(V)dimVm`$, the result follows. Thus, if we are searching for equidimensional quotient morphisms $`\pi _{\widehat{G}}`$, then we must assume that
$`()`$ $`dimN(V)=dimVdimV//G`$ and $`Y_m(N(V))\text{}`$.
In this setting, analogs of results (8)–(8) are proved in a quite similar fashion. Let $`N^u(\widehat{V})`$ denote the zero-fibre of the morphism $`\overline{\pi }:\widehat{V}\text{𝕜}^m`$ defined by
$`\overline{\pi }(v_1,v_2)=(\widehat{F}_1(v_1,v_2),\mathrm{},\widehat{F}_m(v_1,v_2))=((\text{d}f_1)_{v_1},v_2,\mathrm{},(\text{d}f_m)_{v_1},v_2)`$.
8.10 Theorem. Under the assumptions (1)–(3) of Theorem 7 and $`()`$, we have
* The morphism $`\pi _{\widehat{G}}:\widehat{V}\widehat{V}//\widehat{G}`$ (resp. $`\overline{\pi }:\widehat{V}\text{𝕜}^m`$) is equidimensional if and only if $`\mathrm{codim}_{N(V)}Y_i(N(V))mi`$ (resp. $`\mathrm{codim}_VY_imi`$) for all $`i`$.
* If $`N(\widehat{V})`$ (resp. $`N^u(\widehat{V})`$) is irreducible, then $`\pi _{\widehat{G}}`$ (resp. $`\overline{\pi }`$) is equidimensional;
* $`N(\widehat{V})`$ (resp. $`N^u(\widehat{V})`$) is irreducible if and only if $`\mathrm{codim}_{N(V)}Y_i(N(V))mi+1`$ (resp. $`\mathrm{codim}_VY_imi+1`$) for $`i<m`$.
* If $`\mathrm{codim}_{N(V)}Y_i(N(V))a`$, then $`\mathrm{codim}_VY_imi+a`$.
Proof. The proof of parts 1–3 is similar to the proof of Theorem 8. For the last part, we notice that $`dim\overline{Y}_i//Gi`$. Therefore $`dimY_ii+dimY_i(N(V))i+dimN(V)a=dimV(mi+a)`$. (Cf. the proof of Prop. 2.1 in .) Consider the homomorphism of $`\text{𝕜}[V]`$-modules
$$\widehat{\mu }:\mathrm{Mor}(V,V^{})\mathrm{Mor}(V,g^{})$$
defined by $`\widehat{\mu }(F)(v),s:=F(v),sv`$ for $`vV,sg`$. Here “$``$” refers to the $`g`$-module structure on $`V`$ and the first (resp. second) $`,`$ stands for the pairing of $`g`$ and $`g^{}`$ (resp. $`V`$ and $`V^{}`$). By \[27, theorem 1.9\], $`\mathrm{ker}\widehat{\mu }`$ is a free $`\text{𝕜}[V]`$-module of rank $`m`$ generated by $`\text{d}f_i`$, $`i=1,\mathrm{},m`$. Let $`\widehat{E}`$ denote the $`\text{𝕜}[V]`$-module $`\mathrm{Im}\widehat{\mu }`$.
8.11 Theorem. Suppose $`N(\widehat{V})`$ is irreducible. Then
* The condition $`(_2)`$ is satisfied by $`\widehat{E}`$.
* $`\text{Sym}_{\text{𝕜}[V]}(\widehat{E})`$ is a factorial domain of Krull dimension $`2dimVm`$.
* $`N^u(\widehat{V})`$ is an irreducible factorial complete intersection, and $`\text{𝕜}[N^u(\widehat{V})]=\text{Sym}_{\text{𝕜}[V]}\widehat{E}`$.
* $`N^u(\widehat{V})=\overline{\mathrm{Im}(\varkappa )}`$, where $`\varkappa :V\times gV\times V`$ is defined by $`\varkappa (v,x)=(v,xv)`$, $`xg,vV`$.
The proof of Theorem 8 is omitted, since it is similar to the proof of Theorem 8.
## 9. Isotropy contractions and $`_2`$-contractions of semisimple Lie algebras
Let $`h`$ be a subalgebra of $`q`$ such that $`q=h𝔪`$ for some $`\mathrm{ad}h`$-stable subspace $`𝔪q`$. (Such an $`h`$ is said to be reductive in $`q`$.) For instance, if $`\vartheta `$ is an involutory automorphism of $`q`$, then $`+1`$ and $`1`$-eigenspaces of $`\vartheta `$ yield such a decomposition. The fixed-point subalgebra of an involutory automorphism is called a symmetric subalgebra.
9.1 Definition. If $`h`$ is reductive in $`q`$, then the semi-direct product $`h𝔪`$ is called an isotropy contraction of $`q`$. If $`h`$ is symmetric, so the decomposition $`q=h𝔪`$ is a $`_2`$-grading, then $`h𝔪`$ is also called a $`_2`$-contraction of $`q`$. Notice that $`h𝔪`$ is a contraction of $`q`$ in the sense of the deformation theory of Lie algebras, see e.g. \[45, Chapter 7, § 2\]. More precisely, consider the invertible linear map $`c_t:qq`$, $`t\text{𝕜}\{0\}`$, such that $`c_t(h+m)=h+t^1m`$. Define the new Lie algebra multiplication $`[,]_{(t)}`$ on the vector space $`q`$ by the rule
$$[x,y]_{(t)}:=c_t\left([c_t^1(x),c_t^1(y)]\right),x,yq.$$
Then, for all $`t0`$, the algebras $`q_{(t)}`$ are isomorphic, and $`lim_{t0}q_{(t)}=h𝔪`$.
9.2 Lemma. Any Takiff Lie algebra is a $`_2`$-contraction. Proof. Consider the direct sum of Lie algebras $`qq`$ and the involution $`\vartheta `$ permuting the summands. Then the corresponding $`_2`$-contraction is isomorphic to $`qq`$. In the rest of the section, we only consider isotropy contractions such that the initial ambient Lie algebra is semisimple and the subalgebra is reductive. Let $`k=h𝔪`$ be an isotropy contraction of a semisimple Lie algebra $`g`$. For $`g`$, one has equalities
$$\mathrm{rk}g=\mathrm{ind}g=dimg//G.$$
The first natural question is:
To which extent this remains true for isotropy contractions?
Recall that the complexity of a homogeneous space $`G/H`$, denoted $`c(G/H)`$, equals $`\mathrm{trdeg}\text{𝕜}(G/H)^B`$, where $`B`$ is a Borel subgroup of $`G`$, and $`G/H`$ is said to be spherical if $`c(G/H)=0`$. We refer to for basic facts on complexity.
9.3 Proposition.
(1) We have $`\mathrm{ind}k=\mathrm{ind}g+2c(G/H)`$. In particular, $`\mathrm{ind}k=\mathrm{ind}g`$ if and only if $`H`$ is a spherical subgroup of $`G`$.
(2) $`dimk//K=dimz_g(x)`$, where $`xh`$ is an $`h`$-regular semisimple element. Proof. (1) By , $`\mathrm{ind}k=\mathrm{trdeg}\text{𝕜}(m^{})^H+\mathrm{ind}s`$, where $`s`$ is a generic stabiliser for $`(H:m^{})`$. Since $`m`$ is an orthogonal $`h`$-module, there is no difference between $`m`$ and $`m^{}`$, the action $`(H:m)`$ is stable and therefore $`s`$ is reductive. Hence $`\mathrm{ind}k=dimm//H+\mathrm{rk}s`$. On the other hand, there is a formula for $`c(G/H)`$ in terms of the isotropy representation $`(H:m)`$. Namely, $`2c(G/H)=dimm//H\mathrm{rk}g+\mathrm{rk}s`$ \[28, Cor. 2.2.9\]. Hence the conclusion.
(2) By Theorem 6, $`dimk//K=\mathrm{rk}h+dimm^{t_h}`$. The latter equals $`dimz_h(x)+dimm^x`$ for a regular semisimple element $`xt_hh`$. Remark. It is a general fact that the index of a Lie algebra cannot decrease under contraction. The previous result gives a precise meaning for this in case of isotropy contractions.
9.4 Corollary. If $`g=g_0g_1`$ is a $`_2`$-grading and $`k=g_0g_1`$ is the respective $`_2`$-contraction, then $`\mathrm{ind}k=dimk//K=\mathrm{rk}g`$. Proof. As is well known, any symmetric subgroup $`G_0G`$ is spherical, and $`g_0`$ contains a regular semisimple element of $`g`$. Thus, for $`_2`$-contractions one obtains two, usually different, decompositions of the rank of $`g`$:
$$\mathrm{rk}g=\{\begin{array}{c}\mathrm{ind}k=\mathrm{rk}s+dimg_1//G_0;\hfill \\ dimk//K=\mathrm{rk}g_0+dim(g_1)^{t_0},\hfill \end{array}$$
where $`t_0`$ is a Cartan subalgebra of $`g_0`$.
If $`h`$ contains a $`g`$-regular semisimple element, then $`\text{𝕜}[g]^G`$ and $`\text{𝕜}[k]^K`$ are graded polynomial algebras of the same Krull dimension. The second natural question is:
Is there a relationship between the degrees of free homogeneous generators (basic invariants) ?
Let $`\text{Deg}(A)`$ denote the multiset of degrees of free generators of a graded polynomial algebra $`A`$. The elements of $`\text{Deg}(A)`$ are assumed to be increasingly ordered.
9.5 Theorem. (1) If $`h`$ contains a $`g`$-regular semisimple element, then $`\text{Deg}(\text{𝕜}[k]^K)\text{Deg}(\text{𝕜}[g]^G)`$ (componentwise inequalities).
(2) Suppose a regular nilpotent element of $`h`$ is also regular in $`g`$. Then
$`\text{g-exp}_g(g)=\text{g-exp}_h(h)\text{g-exp}_h(m)`$ (the union of multisets). Equivalently, $`\text{Deg}(\text{𝕜}[k]^K)=\text{Deg}(\text{𝕜}[g]^G)`$. Proof. (1) Recall that $`k=lim_{t0}g_{(t)}`$. It is easily seen that this contraction gives rise to ”a curve in the space of algebras of invariants” and to an embedding $`lim_{t0}\text{𝕜}[g_{(t)}]^{G_{(t)}}\text{𝕜}[k]^K`$. The limit exists, because $`\text{𝕜}[g_{(t)}]^{G_{(t)}}`$ is graded and the (finite) dimension the homogeneous component of a given degree does not depend on $`t`$; so that the limit is taken in a suitable Grassmannian.
(2) Let $`\{e,h,f\}`$ be a principal $`sl_2`$-triple in $`h`$ (see \[45, Ch. 6,§ 2.3\]). By the assumption, it is also a principal $`sl_2`$-triple in $`g`$. By a result of R. Brylinski , the generalised exponents of a $`G`$-module $`V`$ are obtained as follows. Take the subspace $`V^T`$ and its “$`e`$-limit” $`lim_e(V^T)V`$, see \[8, § 2\] for the precise definition. Then $`\text{g-exp}_g(V)`$ is the multiset of $`h`$-eigenvalues on $`lim_e(V^T)`$. It is important that this “$`e`$-limit” depends only on the $`\{e,h\}`$-module structure on $`V`$. In our setting, $`g`$ and $`k`$ are isomorphic as $`h`$-modules, and $`k=hm`$ as $`h`$-module. Therefore
$$\text{g-exp}_g(g)=\text{g-exp}_h(g)=\text{g-exp}_h(k)=\text{g-exp}_h(h)\text{g-exp}_h(m).$$
The second assertion follows from Theorem 6, because $`\text{Deg}(\text{𝕜}[k]^K)=\text{g-exp}_h(k)+1`$ (componentwise) and likewise for $`\text{𝕜}[g]^G`$. Part (2) of this theorem can be used for finding generalised exponents of certain representations.
9.6 Example. Let $`g`$ be $`so_8`$ and $`h`$ the exceptional Lie algebra of type $`\text{G}_2`$ ($`dimg=28`$, $`dimh=14`$). The restriction of the defining representation of $`g`$ to $`h`$ is the sum of $`V(7)`$, the 7-dimensional simple $`h`$-module, and a 1-dimensional trivial module. Let $`eh`$ be a regular nilpotent element. It is known that $`V(7)`$ is a cyclic $`e`$-module. Therefore, as element of $`so_8`$, $`e`$ has the Jordan form with blocks of size 7 and 1. Hence $`e`$ is also regular in $`so_8`$. Here $`m=V(7)V(7)`$. Since $`\text{g-exp}_g(g)=\{1,3,3,5\}`$ and $`\text{g-exp}_h(h)=\{1,5\}`$, we conclude that $`\text{g-exp}_h(V(7))=\{3\}`$. That is, the $`\text{𝕜}[h]^H`$-module $`\mathrm{Mor}_h(h,V(7))`$ is generated by the covariant of degree 3.
This is also an instructive illustration to Theorem 8 and Corollary 8. Here $`m=\mathrm{rk}g\mathrm{rk}h=2`$, hence $`\pi _K`$ is equidimensional. The basic covariant in $`\mathrm{Mor}_h(h,V(7))`$ vanishes on the subregular nilpotent orbit in $`𝒩(h)`$. This follows from a result of Broer on the ideal defining the closure of the subregular nilpotent orbit \[7, § 4\]. Therefore $`\mathrm{codim}_{𝒩(h)}X_0(𝒩(h))=2`$ and $`N(k)`$ appears to be reducible. From now on, we assume that $`k`$ is a $`_2`$-contraction of $`g`$.
9.7 Theorem. Let $`g=g_0g_1`$ be a $`_2`$-graded semisimple Lie algebra and $`k=g_0g_1`$ its $`_2`$-contraction. Then $`N(k)`$ is irreducible. Proof. Let $`\vartheta `$ be the involution of $`g`$ determining the $`_2`$-grading. It suffices to handle the case in which $`g`$ is not a sum of $`\vartheta `$-stable ideals. This means that either $`g`$ is simple or $`g=ss`$, where $`s`$ is simple and $`\vartheta `$ permutes the factors. In the second case, $`k=ss`$ is a Takiff Lie algebra, and the required result is proved in \[16, Theorem 2.4\]. Therefore we concentrate on the first case.
From now on, $`g`$ is simple. Write $`𝒩_0`$ for the null-cone in $`g_0`$ and $`K`$ for the Takiff group $`G_0g_1`$. Since $`g_1`$ is an orthogonal $`G_0`$-module, we do not distinguish $`g_1`$ and $`(g_1)^{}`$.
1) Suppose $`\vartheta `$ is inner. Then $`\mathrm{rk}g=\mathrm{rk}g_0`$ and therefore the $`g_0`$-module $`g_1`$ has no zero weight space. As is noted in Section 8, the null-cone $`N(k)`$ is then isomorphic to $`𝒩_0\times g_1`$.
2) Suppose $`\vartheta `$ is outer. This is the difficult part of the proof, which relies on the classification of the involutions of simple Lie algebras. Recall that $`m=dim(g_1)^{t_0}=\mathrm{rk}g\mathrm{rk}g_0`$.
$`(a_1)`$ $`\mathrm{rk}g_0=\mathrm{rk}g1`$ and $`m=1`$. Here the assertion follows from Corollary 8. This happens if $`g=so_{2n}`$ and $`g_0=so_{2k+1}\times so_{2l+1}`$ with $`k+l=n1`$.
$`(a_2)`$ $`\mathrm{rk}g_0=\mathrm{rk}g2`$ and $`m=2`$. By Corollary 8, $`\pi _K`$ is equidimensional. Still, $`N(k)`$ can be reducible a priori. To prove that this is not the case, consider the hierarchy $`X_0(𝒩_0)\overline{X_1(𝒩_0)}\overline{X_2(𝒩_0)}=𝒩_0`$ determined by the basic covariants of type $`g_1`$. Invoking the criterion of irreducibility (Theorem 8(iii)) with $`m=2`$ shows that only the condition with $`i=0`$ has to be satisfied. That is, we must have $`\mathrm{codim}_{𝒩_0}X_0(𝒩_0)3`$. This means that each nilpotent orbit in $`𝒩_0`$ of codimension 2 does not belong to $`X_0(𝒩_0)`$, i.e., there should exist a covariant $`F\mathrm{Mor}_{G_0}(g_0,g_1)`$ that does not vanish on such an orbit.
There are two involutions with $`m=2`$ in the exceptional algebras. In both cases, $`g`$ is of type $`\text{E}_6`$ and $`g_0`$ is either $`\text{F}_4`$ or $`\text{C}_4`$. Furthermore, the degrees of basic covariants of type $`g_1`$ are $`4,\mathrm{\hspace{0.17em}8}`$ in both cases. Since $`g_0`$ is simple here, $`𝒩_0`$ has a unique orbit of codimension 2, the so-called subregular nilpotent orbit $`𝒪_{sub}`$. The closure of $`𝒪_{sub}`$ is normal and the equations of $`\overline{𝒪}_{sub}`$ in $`\text{𝕜}[𝒩_0]`$ are explicitly described, see \[7, § 4\]. Therefore, it is not hard to verify that the covariant of degree 4 survives on $`𝒪_{sub}`$.
$`(a_3)`$ It remains to handle two series of $`(g,g_0)`$: $`(sl_{2n},sp_{2n})`$ and $`(sl_n,so_n)`$. In these cases, we explicitly describe the covariants of type $`g_1`$ and verify that the condition of Theorem 8(iii) is satisfied. Actually, we show that, for all $`_2`$-contractions of simple Lie algebras, a stronger inequality holds, see Eq. (9.9) below. Let us adapt Theorem 8 to our setting. We consider the stratification of $`𝒩_0`$ determined by covariants of type $`g_1`$. Since $`𝒩_0`$ consists of finitely many $`G_0`$-orbits, condition 8(iii) can be verified for each orbit separately. Therefore, it can be written as
(9.8)
$$dimz_{g_0}(x)\mathrm{rk}g_0>mdim\text{span}\{F_1(x),\mathrm{},F_m(x)\}\text{ if }x𝒩_0X_m(𝒩_0),$$
cf. Remark 8(1). What we are going to prove is:
(9.9)
$$dimz_{g_0}(x)\mathrm{rk}g_02\left(mdim\text{span}\{F_1(x),\mathrm{},F_m(x)\}\right)\text{ for any }x𝒩_0.$$
Clearly, the last version is stronger and has an advantage of being stated more uniformly.
9.10 Theorem. Inequality (9.9) holds for any $`_2`$-grading of a simple Lie algebra $`g`$. Proof. Since the difference in the left-hand side of (9.9) is always even, there is no distinction between inequalities (9.8) and (9.9) for $`m2`$. Therefore the proof of Theorem 9 shows that it remains to verify Eq. (9.9) for the following series of $`_2`$-gradings:
* $`g_0=sp(V)`$, $`g_1=_0^2(V)`$, $`dimV=2n`$.
* $`g_0=so(V)`$, $`g_1=𝒮_0^2(V)`$. Here one actually has two series, depending on the parity of $`dimV`$.
We use familiar matrix models of classical Lie algebras and their representations. In the following computations, we need the fact that the nilpotent $`G_0`$-orbits are classified by certain partitions of $`dimV`$, see \[39, IV.2.15\], \[45, Ch. 6 §2.2\]. A minor unpleasant phenomenon related to $`so_{2n}`$ is that there are two isomorphic $`SO_{2n}`$-orbits corresponding to a “very even partition”. This does not affect, however, our computations. For $`x𝒩_0`$, let $`𝜼=(\eta _1,\eta _2,\mathrm{})`$ denote the corresponding partition. Write $`(\widehat{\eta }_1,\widehat{\eta }_2,\mathrm{},\widehat{\eta }_s)`$ for the dual partition. This means in particular that $`s=\eta _1`$. What we need from these partitions is an explicit formula for $`dimz_{g_0}(x)`$ and a way to determine $`i`$ such that $`xX_i(𝒩_0)`$.
Let us begin with the symplectic case. Let $`J`$ be a skew-symmetric non-degenerate bilinear form on $`V`$, which is identified with its matrix in a certain basis for $`V`$. Then
$``$ $`sp_{2n}=sp(V)=sp(V,J)`$ is the space of matrices $`\{xgl(V)xJ\text{ is symmetric}\}`$;
$``$ the representation space $`_0^2(V)`$ can be regarded as the space of skew-symmetric matrices modulo one-dimensional subspace generated by $`J`$. The $`sp_{2n}`$-action on the space of skew-symmetric matrices is given by $`(x,A)xJA+A(xJ)^t`$.
In this case $`m=n1`$, i.e., there are $`n1`$ basic covariants of type $`g_1`$. Since any regular nilpotent element in $`sp_{2n}`$ is also regular in $`sl_{2n}`$, the generalised exponents of the $`g_0`$-module $`g_1`$ can be found using Theorem 9(2). These are $`2,4,\mathrm{},2n2`$. The key observation is that the corresponding covariants have a very simple expression. Namely, consider the maps $`(xsp_{2n})F_i(x)=x^{2i}J`$, $`i=1,\mathrm{},n1`$. It is easily seen that $`x^{2i}J`$ is skew-symmetric and each $`F_i`$ is $`Sp_{2n}`$-equivariant. Because the $`F_i`$’s are linearly independent over $`\text{𝕜}[g_0]^{G_0}`$, these are precisely the basic covariants.
9.11 Proposition. Inequality (9.9) is satisfied for $`(Sp(V),_0^2(V))`$. Proof. By \[17, Corollary 3.8(a)\], the dimension of the centraliser of $`x`$ in $`g_0=sp_{2n}`$ is given by the formula $`dimz_{g_0}(x)={\displaystyle \frac{1}{2}}({\displaystyle \underset{i}{}}\widehat{\eta }_i^2+\mathrm{\#}\{j\eta _j\text{ is odd}\})`$.
The maximal nonzero power of $`x`$ is determined by the size of the maximal Jordan block, i.e., $`\eta _1`$. Therefore $`xX_i(𝒩_0)`$ if and only if $`x^{2i}0`$ and $`x^{2i+2}=0`$ if and only if $`[\frac{\eta _11}{2}]=i`$. Hence inequality (9.9), which we wish to prove, can be written as
$$2\left[\frac{\eta _11}{2}\right]+\frac{1}{2}\left(\underset{i=1}{\overset{s}{}}\widehat{\eta }_i^2+\mathrm{\#}\{j\eta _j\text{ is odd}\}\right)n2(n1)0.$$
Using the relations $`\widehat{\eta }_i=2n`$ and $`\eta _1=s`$, the left-hand side is transformed as follows:
$$\begin{array}{c}2\left[\frac{\eta _11}{2}\right]+\frac{1}{2}\left(\underset{i=1}{\overset{s}{}}\widehat{\eta }_i^2+\mathrm{\#}\{j\eta _j\text{ is odd}\}\right)\frac{3}{2}\widehat{\eta }_i+2=\hfill \\ \hfill 2\left[\frac{s1}{2}\right]+\frac{1}{2}\left(\underset{i=1}{\overset{s}{}}(\widehat{\eta }_i^23\widehat{\eta }_i)+\mathrm{\#}\{j\eta _j\text{ is odd}\}\right)+2=\\ \hfill \frac{1}{2}\left(\underset{i=1}{\overset{s}{}}(\widehat{\eta }_i1)(\widehat{\eta }_i2)+\mathrm{\#}\{j\eta _j\text{ is odd}\}\right)+2\left[\frac{s+1}{2}\right]s.\end{array}$$
The first group of summands is non-negative, and so is the last group. Thus, inequality (9.9) holds for any nilpotent orbit in $`sp_{2n}`$. We continue with the orthogonal case, with $`dimV=N`$. Here $`g_0`$ is the space of skew-symmetric $`N\times N`$-matrices and $`g_1=𝒮_0^2(V)`$ is the space of traceless symmetric $`N\times N`$-matrices.
If $`N=2n+1`$, then $`m=n`$. In this case, a regular nilpotent element of $`so_{2n+1}`$ is also regular in $`sl_{2n+1}`$, so that Theorem 9(2) applies, and $`\text{g-exp}_{G_0}(g_1)=\{2,4,\mathrm{},2n\}`$. Similarly to the symplectic case, we find that $`xF_i(x)=x^{2i}`$, $`i=1,2,\mathrm{},n`$, are the basic covariants. If $`N=2n`$, then $`m=n1`$. A regular nilpotent element of $`so_{2n}`$ is not regular in $`sl_{2n}`$, but $`F_1,\mathrm{},F_{n1}`$ are still the basic covariants. For, the $`F_i`$’s are linearly independent over $`\text{𝕜}[g_1]^{G_0}`$ and neither of them vanishes on the regular nilpotent orbit in $`so_{2n}`$.
9.12 Proposition. Inequality (9.9) is satisfied for $`(SO(V),𝒮_0^2(V))`$. Proof. By \[17, Corollary 3.8(a)\], the dimension of the centraliser of $`x`$ in $`g_0=so_N`$ is given by the formula $`dimz_{g_0}(x)={\displaystyle \frac{1}{2}}({\displaystyle \underset{i}{}}\widehat{\eta }_i^2\mathrm{\#}\{j\eta _j\text{ is odd}\})`$. The constraints imposed on partitions in the orthogonal case imply that $`\widehat{\eta }_1dimV(mod2)`$.
The maximal nonzero power of $`x`$ is determined by the size of the maximal Jordan block. Therefore $`xX_i(𝒩_0)`$ if and only if $`x^{2i}0`$ and $`x^{2i+2}=0`$ if and only if $`[\frac{\eta _11}{2}]=i`$. The following computations are slightly different for $`so_{2n+1}`$ and $`so_{2n}`$.
1. $`N=2n+1`$. Here inequality (9.9) can be written as
$$2\left[\frac{\eta _11}{2}\right]+\frac{1}{2}\left(\underset{i}{}\widehat{\eta }_i^2\mathrm{\#}\{j\eta _j\text{ is odd}\}\right)3n0.$$
Using the relations $`\widehat{\eta }_i=2n+1`$ and $`\eta _1=s`$, the left-hand side is transformed as follows:
$$\begin{array}{c}2\left[\frac{s1}{2}\right]+\frac{1}{2}\left(\underset{i=1}{\overset{s}{}}\widehat{\eta }_i^2\mathrm{\#}\{j\eta _j\text{ is odd}\}\right)\frac{3}{2}(\underset{i=1}{\overset{s}{}}\widehat{\eta }_i1)=\hfill \\ \hfill 2\left[\frac{s1}{2}\right]+\frac{1}{2}\left(\underset{i=1}{\overset{s}{}}(\widehat{\eta }_i^23\widehat{\eta }_i+2)2s+3\mathrm{\#}\{j\eta _j\text{ is odd}\}\right)=\\ \hfill \frac{1}{2}(\underset{i=1}{\overset{s}{}}(\widehat{\eta }_i1)(\widehat{\eta }_i2)\mathrm{\#}\{j\eta _j\text{ is odd}\}+4\left[\frac{s+1}{2}\right]2s1)=:𝖫.\end{array}$$
To see that $`𝖫`$ is nonnegative, consider several cases.
(a) $`\widehat{\eta }_1=1`$ and hence all $`\widehat{\eta }_i=1`$. Then $`s=2n+1`$ and $`𝖫=0`$.
(b) $`\widehat{\eta }_1=3`$ and therefore $`𝜼=(\eta _1,\eta _2,\eta _3)`$. Then $`_{i=1}^s(\widehat{\eta }_i1)(\widehat{\eta }_i2)=2\eta _3`$. Hence
$$𝖫=\eta _3+2\left[\frac{\eta _1+1}{2}\right]\eta _1\frac{1}{2}(1+\mathrm{\#}\{j\eta _j\text{ is odd}\}).$$
Taking into account that the even parts in $`(\eta _1,\eta _2,\eta _3)`$ occur pairwise and $`\eta _1+\eta _2+\eta _3`$ is odd, one quickly verifies that $`𝖫`$ is always nonnegative.
(c) $`\widehat{\eta }_15`$. Then $`_{i=1}^s(\widehat{\eta }_i1)(\widehat{\eta }_i2)\widehat{\eta }_1+2\mathrm{\#}\{j\eta _j\text{ is odd}\}+2`$. Next,
$`4\left[\frac{s+1}{2}\right]2s11`$. Hence $`𝖫`$ is positive.
Thus, inequality (9.9) holds for any nilpotent orbit in $`so_{2n+1}`$.
2. $`N=2n`$. Here the inequality we need to prove reads
$$2\left[\frac{\eta _11}{2}\right]+\frac{1}{2}\left(\underset{i}{}\widehat{\eta }_i^2\mathrm{\#}\{j\eta _j\text{ is odd}\}\right)n2(n1)0.$$
Using the relations $`\widehat{\eta }_i=2n`$ and $`\eta _1=s`$, the left-hand side is being transformed to
$$\frac{1}{2}(\underset{i=1}{\overset{s}{}}(\widehat{\eta }_i1)(\widehat{\eta }_i2)\mathrm{\#}\{j\eta _j\text{ is odd}\})+2\left[\frac{s+1}{2}\right]s=:𝖫.$$
Again, consider some cases.
(a) $`\widehat{\eta }_1=2`$, i.e., $`x`$ has only two Jordan blocks $`(\eta _1,\eta _2)`$. Then $`\eta _1,\eta _2`$ have the same parity, and in both cases $`𝖫=0`$.
(b) $`\widehat{\eta }_14`$. Then $`_{i=1}^s(\widehat{\eta }_i1)(\widehat{\eta }_i2)\widehat{\eta }_1+2>\mathrm{\#}\{j\eta _j\text{ is odd}\}`$. Since $`2\left[\frac{s+1}{2}\right]s0`$, the total expression is positive.
Thus, inequality (9.9) holds for any nilpotent orbit in $`so_{2n}`$. This completes the proof of Theorem 9. Thus, all verifications needed to complete the proof of Theorem 9 are done. Below, we gather our results on $`_2`$-contractions of semisimple Lie algebras.
9.13 Theorem. Let $`k=g_0g_1`$ be a $`_2`$-contraction of a semisimple Lie algebra $`g`$. Then
* $`\text{𝕜}[k]^K`$ is a polynomial algebra of Krull dimension $`\mathrm{rk}g`$;
* $`N(k)`$ is an irreducible complete intersection. If $`\text{𝕜}[k]^K=\text{𝕜}[f_1,\mathrm{},f_l]`$, $`l=\mathrm{rk}g`$, then the ideal of $`N(k)`$ in $`\text{𝕜}[k]`$ is generated by $`f_1,\mathrm{},f_l`$.
* the quotient morphism $`\pi _K:kk//K`$ is equidimensional;
* $`\text{𝕜}[k]`$ is a free $`\text{𝕜}[k]^K`$-module.
* if $`\kappa :g_0g_1g_0g_1`$ is defined by $`\kappa (x_0,x_1)=(x_0,[x_0,x_1])`$, then $`\overline{\mathrm{Im}\kappa }=N^u(q)`$ and it is a factorial complete intersection of codimension $`\mathrm{rk}g\mathrm{rk}g_0`$.
* the coadjoint representation of $`k`$ has a generic stabiliser.
Proof. Part (1) follows from Theorem 6. The irreducibility in Part (2) is just Theorem 9. Let $`x𝒩_0`$ be a regular nilpotent element. Then $`\stackrel{~}{x}=(x,0)N(k)`$, and the description of basic invariants $`f_1,\mathrm{},f_l`$ in Theorem 6 shows that $`(\text{d}f_1)_{\stackrel{~}{x}},\mathrm{},(\text{d}f_l)_{\stackrel{~}{x}}`$ are linearly independent. Then a standard argument shows that the ideal of $`N(k)`$ is generated by $`f_1,\mathrm{},f_l`$ (cf. \[21, Prop. 6\].) Part (3) follows from (2) and Theorem 8(2). Part (4) is a formal consequence of Parts (1) and (3). Part (5) follows from Theorem 8 and the irreducibility of $`N(k)`$. Since the isotropy representation of any symmetric subalgebra of $`g`$ is polar, part (6) follows from Theorem 5. To prove the irreducibility of $`N(k)`$, inequality (9.8) is sufficient. However, our efforts in proving stronger inequality (9.9) are not in vain, because that result also has a geometric meaning.
9.14 Theorem. Let $`g=g_0g_1`$ be a $`_2`$-grading. Consider the semi-direct product $`\stackrel{~}{k}=g_0(g_1g_1)`$ and the corresponding adjoint representation $`(\stackrel{~}{K}:\stackrel{~}{k})`$. Then the quotient morphism $`\pi _{\stackrel{~}{K}}`$ is equidimensional. Proof. The criterion for equidimensionality of $`\pi _{\stackrel{~}{K}}`$, Theorem 8(i), written out in this case yields precisely inequality (9.9). Main efforts in this section were devoted to $`_2`$-contractions of $`g`$. However, there are interesting examples of other isotropy contractions with full bunch of good properties.
9.15 Examples. 1. Suppose $`g=so_7`$ and $`h`$ is a simple subalgebra of type $`\text{G}_2`$. It is a ”truncation” of Example 9. Here $`m=V(7)`$, and one easily verifies that all conclusions of Theorem 9 hold for $`k=hm`$.
2. $`g=sl_{2n+1}`$ and $`h=sp_{2n}=sp(V)`$. Here the $`sp(V)`$-module $`m`$ equals $`^2(V)VV`$. Since the $`sp(V)`$-module $`V`$ has no zero-weight space, the structure of $`N(hm)`$ is essentially the same as for the $`_2`$-contraction of the symmetric pair $`(sl_{2n},sp_{2n})`$. Remark. Our proofs of Theorems 9 and 9 use classification of involutory automorphisms and explicit considerations of cases. It would be extremely interesting to find a case-free proof for the irreducibility of $`N(k)`$. Especially, because the corresponding irreducibility result for the Takiff algebra $`gg`$ can derived without checking cases. We discuss this topic in the following section.
## 10. Reductive Takiff Lie algebras and their representations
The attentive reader may have noticed that we stated and proved the stronger inequality (9.9) only for the $`_2`$-gradings of simple Lie algebras, leaving aside the permutation of two factors in $`g\times g`$ and the corresponding Takiff algebra $`\widehat{g}`$.
The situation here is as follows. By Theorem 8, the counterpart of inequality (9.8) for $`\widehat{g}`$ is equivalent to the irreducibility of $`N(\widehat{g})`$, and this was already proved by Geoffriau . His proof consists of explicit verifications for all simple types. It was noticed by M. Brion that a classification-free proof of (9.8) for $`\widehat{g}`$, and hence the irreducibility of $`N(\widehat{g})`$, can be derived from the fact that $`𝒩`$ is a complete intersection having only rational singularities, see below. The advantage of the Takiff algebra case is that the rather mysterious term $`dim\text{span}\{F_1(x),\mathrm{},F_m(x)\}`$ is being interpreted as the rank of the differential of the quotient map $`\pi _G:gg//G`$ at $`x`$.
On the other hand, we will prove here the counterpart of (9.9) for $`\widehat{g}`$, using the classification. Brion’s idea cannot be applied directly to obtain a case-free proof of that stronger result. The reason for being interested in proving a counterpart of (9.9) for $`\widehat{g}`$ is that we deduce from this the equidimensionality of some other quotient morphisms, see Theorems 10,10.
We work in the setting of case B) from Section 8. 10.1 Definition. Let $`\rho :GGL(V)`$ be a representation of a connected reductive group $`G`$. Then $`V`$ or $`\rho `$ is said to be extremely good if
* $`\text{𝕜}[V]^G`$ is a polynomial algebra;
* $`\mathrm{max}dim_{xV}Gx=dimVdimV//G`$;
* If $`\pi _G:VV//G`$ is the quotient morphism, then $`\mathrm{\Omega }:=\{xV(\text{d}\pi _G)_x\text{ is onto }\}`$ is a big open subset of $`V`$;
* $`N(V):=\pi _G^1(\pi _G(0))`$ consists of finitely many $`G`$-orbits.
* $`N(V)`$ is irreducible and has only rational singularities;
Note that properties (1)–(3) are those appearing in Theorem 7. Recall from Section 7 that if $`G`$ is semisimple, then (2) and (3) are always satisfied.
10.2 Theorem. Let $`V`$ be an extremely good $`G`$-module and $`\widehat{V}=VV`$ the corresponding $`\widehat{G}`$-module. Then
* $`N^{\widehat{G}}(\widehat{V})=N(\widehat{V})`$ is an irreducible complete intersection;
* the ideal of $`N(\widehat{V})`$ in $`\text{𝕜}[\widehat{V}]`$ is generated by the basic invariants in $`\text{𝕜}[\widehat{V}]^{\widehat{G}}`$;
* $`\pi _{\widehat{G}}:\widehat{V}\widehat{V}//\widehat{G}`$ is equidimensional and $`\text{𝕜}[\widehat{V}]`$ is a free $`\text{𝕜}[\widehat{V}]^{\widehat{G}}`$-module.
Proof. Let $`f_1,\mathrm{},f_m`$ be algebraically independent generators of $`\text{𝕜}[V]^G`$. By Theorem 7, $`\text{𝕜}[\widehat{V}]^{\widehat{G}}`$ is freely generated by the polynomials $`f_1,\mathrm{},f_m,\widehat{F}_{f_1},\mathrm{},\widehat{F}_{f_m}`$. Recall from Section 8 the stratification of the null-cone:
$$Y_i(N(V))=\{vN(V)\mathrm{rk}(\text{d}\pi _G)_v=i\},i=0,1,\mathrm{},m.$$
Since $`N(V)`$ contains finitely many $`G`$-orbits, $`\pi _G`$ is equidimensional. If $`Gx`$ is the dense $`G`$-orbit in $`N(V)`$, then $`dimGx=dimVm`$ and therefore $`xY_m(N(V))`$ \[20, Korollar 2\]. (Corollary 2 is stated in Knop’s article under the assumption that $`G`$ is semisimple. However, that proof works also for reductive groups as long as conditions (2) and (3) are satisfied.) Since $`N(V)`$ is irreducible and $`Y_m(N(V))\text{}`$, it is a complete intersection. The condition of the irreducibility of $`N(\widehat{V})`$ (Theorem 8(iii)) can be written as
(10.3)
$$dimVdimGv+\mathrm{rk}(\text{d}\pi _G)_v>2m\text{ if }vY_m(N(V)).$$
We derive this inequality from a property of the local ring of (the closure of) the orbit $`GvN(V)`$. Let $`O`$ be this local ring. Then $`dimO=dimN(V)dimGv`$ and $`\mathrm{edim}O=dimT_vN(V)dimGv=dimg_v\mathrm{rk}(\text{d}\pi _G)_v`$. Here $`\mathrm{edim}O`$ is the embedding dimension of $`O`$ and $`T_vN(V)`$ is the tangent space of $`N(V)`$ at $`v`$. Since $`N(V)`$ has only rational singularities, so has $`O`$. By a result of Goto-Watanabe (see \[25, Theorem 2’\]), if a local ring $`O`$ is a complete intersection with only rational singularities and $`dimO>0`$, then $`\mathrm{edim}O<2dimO`$. Using the above expressions for $`\mathrm{edim}O`$ and $`dimO`$, one obtains inequality (10.3), and thereby the irreducibility of $`N(\widehat{V})`$.
All other statements of the theorem are consequences of the fact that $`N(\widehat{V})`$ is irreducible. By Theorem 8(ii), $`\pi _{\widehat{G}}`$ is equidimensional. If $`vY_m(N(V))`$, then the differentials of the generators $`f_1,\mathrm{},f_m,\widehat{F}_{f_1},\mathrm{},\widehat{F}_{f_m}`$ are linearly independent at $`(v,0)N(\widehat{V})\widehat{V}`$. This fact and the irreducibility of $`N(\widehat{V})`$ imply that $`N(\widehat{V})`$ is a complete intersection whose ideal is generated by the polynomials $`f_1,\mathrm{},f_m,\widehat{F}_{f_1},\mathrm{},\widehat{F}_{f_m}`$ (cf. \[21, Prop. 6\]). Remark. The most subtle point in the definition of extremely good representations is the rationality of singularities of $`N(V)`$. For the adjoint representations, this result is due to W. Hesselink . The idea to exploit the fact that $`𝒩=N(g)`$ is a complete intersection with only rational singularities, and to use the Goto-Watanabe inequality for local rings is due to M. Brion . Since $`(G,\mathrm{Ad})`$ is extremely good, this approach yields a conceptual proof of \[16, Theorem 2.4\].
10.4 Corollary. If $`V`$ is extremely good, then the closure of the image of the map
$$\varkappa :V\times gV\times V,(v,x)(v,xv),$$
is a factorial complete intersection of codimension $`m=dimV//G`$ and the ideal of $`\overline{\mathrm{Im}\varkappa }`$ is generated by $`\widehat{F}_{f_1},\mathrm{},\widehat{F}_{f_m}`$. Proof. This follows from the irreducibility of $`N(\widehat{V})`$ and Theorem 8. Since conditions (4) and (5) are quite restrictive, there are only a few extremely good representations. Below is a list of such irreducible representations known to this author such that $`G`$ is simple and $`\text{𝕜}[V]^G\text{𝕜}`$, except the adjoint ones:
$`(\text{B}_n\text{ or }\text{D}_n,\phi _1),(\text{B}_3,\phi _3),(\text{B}_4,\phi _4),(\text{G}_2,\phi _1),(\text{A}_n,2\phi _1),(\text{A}_{2n1},\phi _2),(\text{E}_6,\phi _1),`$
$`(\text{C}_3,\phi _3),(\text{A}_5,\phi _3),(\text{D}_6,\phi _6),(\text{E}_7,\phi _1),(\text{B}_5,\phi _5),(\text{F}_4,\phi _1),(\text{C}_n,\phi _2)`$.
The representations are given by their highest weights, and $`\{\phi _i\}`$ are fundamental weights of $`G`$ with numbering from . For all representations in the list but the last one, the algebra of covariants, $`\text{𝕜}[V]^U`$, is polynomial (here $`U`$ is a maximal unipotent subgroup of $`G`$). Therefore the same is true for $`\text{𝕜}[N(V)]^U`$. Then a result of Kraft (see \[3, 1.5-6\]) shows that $`N(V)`$ has rational singularities.
I conjecture that if $`G`$ is simple and $`V`$ is a simple $`G`$-module, then $`V`$ is extremely good if and only if $`dimVdimG`$. Practically, this means that one has to only verify that $`N(V)`$ has rational singularities for the following representations:
$`(\text{A}_6,\phi _3),(\text{A}_7,\phi _3),(\text{B}_6,\phi _6),(\text{D}_7,\phi _7)`$.
For $`V=g`$, inequality (10.3) reads
(10.5)
$$dimz_g(x)+\mathrm{rk}(\text{d}\pi _G)_x>2\mathrm{r}\mathrm{k}g=2m\text{ if }xX_m(𝒩)=𝒩^{reg}.$$
This inequality was proved in \[16, 2.6-2.15\] in a case-by-case fashion. Below, we prove a stronger result, which is the counterpart of inequality (9.9) in the context of Takiff algebras. By the Morozov-Jacobson theorem \[45, Ch. 3, Theorem 1.3\], any $`x𝒩\{0\}`$ can be embedded in an $`sl_2`$-triple $`\{x,h,y\}`$, where $`h`$ is semisimple; $`x`$ is said to be even if the $`\mathrm{ad}h`$-eigenvalues in $`g`$ are even. Following E.B. Dynkin, $`h`$ is called a characteristic of $`x`$.
10.6 Theorem. Let $`g`$ be a simple Lie algebra and $`x𝒩`$. Then
(10.7)
$$𝖫:=dimz_g(x)+2\mathrm{r}\mathrm{k}(\text{d}\pi _G)_x3\mathrm{r}\mathrm{k}g0.$$
If $`g=sl_n`$, then $`𝖫=0`$ if and only if the matrix $`x`$ has at most two Jordan blocks. Furthermore, if $`gsl_{2n+1}`$, then $`𝖫=0`$ if and only if $`x`$ is even and $`[z_g(h),z_g(h)]`$ is a sum of several copies of $`sl_2`$. (Here $`h`$ is a characteristic of $`x`$). Proof. The proof is case-by-case. However, the computations themselves are much shorter and more transparent than those in , because our inequality is stronger, and we use formulae for $`dimz_g(x)`$ in terms of dual partitions (already used for $`Sp`$ and $`SO`$ in the proof of Propositions 9 and 9).
For the classical series, we work with the partition of $`x`$; while for the exceptional algebras the explicit classification of nilpotent orbits is used. If $`g=g(𝕍)`$ is classical and $`xg(𝕍)`$ is nilpotent, then $`𝜼=(\eta _1,\eta _2,\mathrm{})`$ is the partition of $`dim𝕍`$ corresponding to $`x`$ and $`(\widehat{\eta }_1,\mathrm{},\widehat{\eta }_s)`$ is the dual partition. Here $`s=\eta _1`$. For $`Sp`$ and $`SO`$, our analysis is quite similar to those in Propositions 9 and 9.
(A) $`g=sl(𝕍)`$, $`dim𝕍=n+1`$.
Here $`dimz_g(x)=_{i=1}^s\widehat{\eta }_i^21`$ and $`\mathrm{rk}(\text{d}\pi _G)_x=\eta _11=s1`$ \[35, Theorem 4.2.1\]. Then
$$𝖫=\underset{i=1}{\overset{s}{}}\widehat{\eta }_i^21+2(s1)3n=\underset{i=1}{\overset{s}{}}\widehat{\eta }_i^23\underset{i=1}{\overset{s}{}}\widehat{\eta }_i+2s=\underset{i=1}{\overset{s}{}}(\widehat{\eta }_i1)(\widehat{\eta }_i2)0.$$
This expression equals zero if and only if all $`\widehat{\eta }_i2`$, i.e., $`x`$ has at most two Jordan blocks.
(B) $`g=so(𝕍)`$, $`dim𝕍=2n+1`$.
Here $`\widehat{\eta }_1`$ is odd, $`dimz_g(x)=\frac{1}{2}\left(_i\widehat{\eta }_i^2\mathrm{\#}\{j\eta _j\text{ is odd}\}\right)`$, and $`\mathrm{rk}(\text{d}\pi _G)_x=[s/2]`$ \[35, Theorem 4.3.3\]. Then
$$\begin{array}{c}𝖫=\frac{1}{2}\left(\underset{i=1}{\overset{s}{}}\widehat{\eta }_i^2\mathrm{\#}\{j\eta _j\text{ is odd}\}\right)+2[s/2]3n=\hfill \\ \hfill \frac{1}{2}\left(\underset{i=1}{\overset{s}{}}(\widehat{\eta }_i1)(\widehat{\eta }_i2)+3\mathrm{\#}\{j\eta _j\text{ is odd}\}2(s2[s/2])\right).\end{array}$$
If $`\widehat{\eta }_1=1`$, then $`𝖫=0`$. This is the case of regular nilpotent elements.
If $`\widehat{\eta }_1=3`$, then $`_{i=1}^s(\widehat{\eta }_i1)(\widehat{\eta }_i2)=2\eta _32`$. Therefore $`2𝖫=2\eta _2+3\mathrm{\#}\{j\eta _j\text{ is odd}\}2(s2[s/2])`$. Since $`\mathrm{\#}\{j\eta _j\text{ is odd}\}3`$ and $`2(s2[s/2])2`$, $`2𝖫`$ is nonnegative. Furthermore, $`𝖫=0`$ if and only if $`\eta _3=1`$ and all $`\eta _i`$’s are odd.
If $`\widehat{\eta }_15`$, then $`_{i=1}^s(\widehat{\eta }_i1)(\widehat{\eta }_i2)\widehat{\eta }_1+2\mathrm{\#}\{j\eta _j\text{ is odd}\}+2`$. Next,
$`32(s2[s/2])0`$. Hence $`𝖫`$ is positive.
(C) $`g=sp(𝕍)`$, $`dim𝕍=2n`$.
Here $`dimz_g(x)=\frac{1}{2}\left(_i\widehat{\eta }_i^2+\mathrm{\#}\{j\eta _j\text{ is odd}\}\right)`$ and $`\mathrm{rk}(\text{d}\pi _G)_x=[s/2]`$ \[35, Theorem 4.3.3\]. Then
$$\begin{array}{c}𝖫=\frac{1}{2}\left(\underset{i=1}{\overset{s}{}}\widehat{\eta }_i^2+\mathrm{\#}\{j\eta _j\text{ is odd}\}\right)+2[s/2]3n=\hfill \\ \hfill \frac{1}{2}\left(\underset{i=1}{\overset{s}{}}(\widehat{\eta }_i1)(\widehat{\eta }_i2)+\mathrm{\#}\{j\eta _j\text{ is odd}\}2(s2[s/2])\right).\end{array}$$
It is easily seen that $`𝖫=0`$ if and only if $`\widehat{\eta }_12`$. Otherwise it is positive.
(D) $`g=so(𝕍)`$, $`dim𝕍=2n`$.
Here $`\widehat{\eta }_1`$ is even and $`dimz_g(x)`$ is as in (B). For the rank of $`\text{d}\pi _G`$, we have \[35, Theorem 4.4.2\]
$`\mathrm{rk}(\text{d}\pi _G)_x=\{\begin{array}{cc}\hfill [s/2],& \text{ if }\widehat{\eta }_14;\hfill \\ \hfill (2ni+1)/2,& \text{ if }𝜼=(2ni,i)\text{ with }i\text{ odd};\hfill \\ \hfill l,& \text{ if }𝜼=(n,n)\text{ and }n=2l.\hfill \end{array}`$
Then
$$\begin{array}{c}𝖫=\frac{1}{2}\left(\underset{i=1}{\overset{s}{}}\widehat{\eta }_i^2\mathrm{\#}\{j\eta _j\text{ is odd}\}\right)+2\mathrm{r}\mathrm{k}(\text{d}\pi _G)_x3n=\hfill \\ \hfill \frac{1}{2}\left(\underset{i=1}{\overset{s}{}}(\widehat{\eta }_i1)(\widehat{\eta }_i2)+4\mathrm{r}\mathrm{k}(\text{d}\pi _G)_x2s\mathrm{\#}\{j\eta _j\text{ is odd}\}\right).\end{array}$$
Now, a consideration of cases shows that $`𝖫=0`$ if $`\widehat{\eta }_1=2`$. If $`\widehat{\eta }_14`$, then $`𝖫>0`$ unless $`𝜼=(\eta _1,\eta _2,1,1)`$, where $`\eta _1,\eta _2`$ are both odd.
(EFG) $`g`$ is exceptional.
It is enough to check inequality (10.7) for sufficiently large orbits (with $`dimz_g(x)3\mathrm{r}\mathrm{k}g`$). To this end, one can consult the tables in \[10, Ch. 8\] for dimensions of orbits and \[35, Appendix\] for the values of $`\mathrm{rk}(\text{d}\pi _G)_x`$. Below we list all non-regular nilpotent orbits with $`𝖫=0`$. The orbits are represented by their Dynkin-Bala-Carter labels.
* $`\text{G}_2(a_1)`$;
* $`\text{F}_4(a_1)`$, $`\text{F}_4(a_2)`$;
* $`\text{E}_6(a_1)`$, $`\text{D}_5`$, $`\text{E}_6(a_3)`$;
* $`\text{E}_7(a_1)`$, $`\text{E}_7(a_2)`$, $`\text{E}_6`$, $`\text{E}_6(a_1)`$;
* $`\text{E}_8(a_1)`$, $`\text{E}_8(a_2)`$, $`\text{E}_8(a_3)`$, $`\text{E}_8(a_4)`$.
Inspecting the tables in \[10, Ch. 8\] shows that these are precisely the even nilpotent orbits whose weighted Dynkin diagrams have no adjacent zeros, which exactly means that $`[z_g(h),z_g(h)]`$ is a sum of several $`sl_2`$’s.
For $`g`$ classical, there is a rule for writing out the characteristic $`h`$ in terms of $`𝜼`$ \[39, Ch. IV\]. Hence the Levi subalgebra $`z_g(h)`$ can be computed. This yields the last assertion of the theorem. A geometric meaning of (10.7) will be made clear in the following result. Let $`g=g_0g_1`$ be a $`_2`$-grading and $`\vartheta `$ the corresponding involutory automorphism of $`g`$. Then $`\vartheta `$ can be extended to an involution of the Takiff algebra $`\widehat{g}`$ by letting $`\vartheta (x+ϵy)=\vartheta (x)+ϵ\vartheta (y)`$. The corresponding eigenspaces are $`\widehat{g}_0=g_0g_0`$ and $`\widehat{g}_1=g_1g_1`$. Here $`\widehat{g}_1`$ is a $`\widehat{g}_0`$-module just in the sense of definition given in Section 7. The $`G_0`$-module $`g_1`$ is not extremely good, so that Theorem 10 cannot be applied. But it is ‘good enough’ in the sense that it satisfies properties (1), (2), (4) of Definition 10.
10.8 Theorem. Suppose $`g=g_0g_1`$ is a $`_2`$-grading of maximal rank, i.e., $`g_1`$ contains a Cartan subalgebra of $`g`$. Then the quotient morphism $`\widehat{\pi }:\widehat{g}_1\widehat{g}_1//\widehat{G}_0`$ is equidimensional. Proof. Recall the relationship between orbits and null-cones for the actions $`(G:g)`$ and $`(G_0:g_1)`$. The null-cones are $`𝒩`$ and $`N(g_1)`$, respectively.
* $`N(g_1)=𝒩g_1`$;
* $`Gxg_1`$ is a union of finitely many $`G_0`$-orbits;
* If $`xg_1`$, then $`dimG_0x=\frac{1}{2}dimGx`$;
* For any $`xg`$, we have $`Gxg_1\text{}`$;
* $`\text{𝕜}[g]^G\text{𝕜}[g_1]^{G_0}`$.
The first three properties hold for all $`_2`$-gradings, whereas the last two are characteristic for the involutions of maximal rank, see .
Let us see what the equidimensionality criterion (Theorem 8(i)) means here. We have $`V=g_1`$, $`G=G_0`$, and $`m=dimg_1//G_0`$. Since $`N(g_1)`$ consists of finitely many $`G_0`$-orbits, that criterion reads
$$dimN(g_1)dimG_0xdimg_1//G_0\mathrm{rk}(\text{d}\pi _{G_0})_x$$
for any $`xN(g_1)`$. Here $`\pi _{G_0}:g_1g_1//G_0`$ is the quotient morphism. In view of the above properties of such $`_2`$-gradings, we have $`dimN(g_1)=\frac{1}{2}dim𝒩=\frac{1}{2}(dimg\mathrm{rk}g)`$, $`dimg_1//G_0=\mathrm{rk}g`$, and $`\mathrm{rk}(\text{d}\pi _{G_0})_x=\mathrm{rk}(\text{d}\pi _G)_x`$. The latter stems from the isomorphism $`\text{𝕜}[g]^G\text{𝕜}[g_1]^{G_0}`$. Rewriting the previous inequality using this data yields precisely inequality (10.7) ! Thus, the fact that $`\widehat{\pi }`$ is equidimensional is essentially equivalent to Theorem 10. Yet another geometric application of Eq. (10.7) is the following (cf. Theorem 9):
10.9 Theorem. Set $`g^{[n]}=gg^n`$, where $`g^n`$ (the sum of $`n`$ copies) regarded as commutative Lie algebra and $`n1`$. Consider the adjoint action $`(G^{[n]}:g^{[n]})`$. Then $`\pi _{G^{[n]}}`$ is equidimensional if and only if $`n2`$. Proof. For $`n=1`$, the assertion is already proved. Next, $`dim(g^n)^T=n\mathrm{rk}g`$ and for $`x𝒩`$ the equidimensionality condition of Theorem 8(i) reads $`dimz_g(x)\mathrm{rk}gn(\mathrm{rk}g\mathrm{rk}(\text{d}\pi _G)_x)`$, which is exactly (10.7) for $`n=2`$. Conversely, if $`n3`$, then this condition is not satisfied for the subregular nilpotent orbit. Remark. In the last theorem, the null-cone $`N(g^{[2]})`$ is always reducible. Indeed, each nilpotent $`G`$-orbit such that $`𝖫=0`$ in (10.7) gives rise to an irreducible component of $`N(g^{[2]})`$, see Remark 8(1). The proof of Theorem 10 shows that, for any $`g`$, there are at least two orbits with $`𝖫=0`$.
There are several equivalent ways to present inequality (10.7). Let $``$ denote the variety of Borel subgroups of $`G`$. For any $`x𝒩`$, set $`_x=\{B^{}x\mathrm{Lie}B^{}\}`$. Recall that $`X_i=X_{i,g}=\{xg\mathrm{rk}(\text{d}\pi _G)_x=i\}`$ and $`X_{i,g}(𝒩)=X_{i,g}𝒩`$. This stratification is determined by the covariants of type $`g`$.
10.10 Proposition. Let $`g`$ be a simple Lie algebra and $`m=\mathrm{rk}g`$. Then the following holds:
(1) $`\mathrm{codim}_𝒩X_{i,g}(𝒩)2(mi)`$ for any $`i=0,1,\mathrm{},m`$;
(2) $`dim_x+\mathrm{rk}(\text{d}\pi _G)_x\mathrm{rk}g`$ for any $`x𝒩`$;
(3) If $`𝒪`$ is the local ring of any $`G`$-orbit in $`𝒩`$, then $`\mathrm{edim}𝒪\frac{3}{2}dim𝒪`$;
(4) If $`g=g_0g_1`$ is a $`_2`$-grading of maximal rank and $`xN(g_1)`$, then $`dim(G_0)_x+\mathrm{rk}(\text{d}\pi _{G_0})_x\mathrm{rk}g`$.
(5) If $`g=g_0g_1`$ is a $`_2`$-grading of maximal rank and $`𝒪^{}`$ is the local ring of a $`G_0`$-orbit in $`N(g_1)`$, then $`\mathrm{edim}𝒪^{}2dim𝒪^{}`$. Proof. In fact, all these conditions are equivalent to inequality (10.7). Since $`𝒩`$ contains finitely many $`G`$-orbits, (1) can be written as $`\mathrm{codim}_𝒩Gx2(m\mathrm{rk}(\text{d}\pi _G)_x)`$ for any $`x𝒩`$, which makes it clear that (1) is equivalent to (10.7). For (2), one should use the fact that $`dim_x=\frac{1}{2}(dimz_g(x)\mathrm{rk}g)`$, see e.g. \[38, 4.3.10, 4.5\]. For (3), one have to use formulae for $`dim𝒪`$ and $`\mathrm{edim}𝒪`$ written out in the proof of Theorem 10. For (4), we notice that since $`dimg_1dimg_0=\mathrm{rk}g`$, the equality $`dimG_0x=\frac{1}{2}dimGx`$ is equivalent to that $`dim(G_0)_x=\frac{1}{2}(dimz_g(x)\mathrm{rk}g)=dim_x`$. Finally, the inequalities in (4) and (5) are obtained from each other via simple transformations. Remark. Concerning (5), we note that this inequality is weaker than the Goto-Watanabe inequality from the proof of Theorem 10, but $`N(g_1)`$ is not normal and can be reducible.
10.11 Corollary. $`\mathrm{codim}_gX_{i,g}3(mi)`$ for any $`i=0,1,\mathrm{},m`$. Proof. It follows from the definition of $`X_{i,g}`$ that $`dim\overline{X}_{i,g}//G=i`$. Since $`\overline{X}_{i,g}`$ is conical, the fibre of the origin of the morphism $`\overline{X}_{i,g}\overline{X}_{i,g}//G`$ has the maximal dimension, i.e., $`dimX_{i,g}i+dimX_{i,g}(𝒩)`$, which is exactly what we want, in view of Proposition 10(1). There are many open problems and observations related to our results on reductive Takiff algebras and $`_2`$-contractions. Here are some of them.
$`1^o`$. It seems that if $`H`$ is a spherical reductive subgroup of $`G`$ and $`k=hm`$ is the corresponding isotropy contraction of $`g`$, then $`\pi _K`$ is always equidimensional. At least, I have verified this in case $`G`$ is simple. In fact, Examples 9 and 9 present several instances of this verification.
$`2^o`$. It would be quite interesting to have a case-free proof for Theorem 10 or, equivalently, 10. Various equivalent forms of that result presented in Proposition 10 suggest that there might be different approaches to proving it. From the geometric point of view, the equidimensionality of $`\widehat{\pi }`$ means that there exists a transversal subspace to $`N(\widehat{g}_1)`$, i.e., a subspace $`U`$ such that $`dimU=dim\widehat{g}_1//\widehat{G}_0`$ and $`UN(\widehat{g}_1)=\{0\}`$.
$`3^o`$. Whenever some quotient morphism is equidimensional, it is interesting to find a natural transversal subspace to the null-cone. One may ask for such a subspace in the setting of Theorems 9, 10, 10. Even for the adjoint representation of $`\widehat{g}=gg`$ it is not known how to naturally construct a transversal space to $`N(\widehat{g})`$. If $`\mathrm{\Delta }_tgg`$ is the diagonally embedded Cartan subalgebra, then $`\mathrm{\Delta }_tN(\widehat{g})=\{0\}`$, so that one has a ”one-half” of a transversal space. The problem is to construct the second half. Similarly, if $`k`$ is a $`_2`$-contraction of a simple Lie algebra, it is not known how to construct a transversal space to $`N^u(k)`$.
$`4^o`$. If one knows that some null-cone $`N`$ is irreducible, then it is tempting to find a resolution of singularities for $`N`$.
$`5^o`$. A case-by-case verification shows that $`X_{1,g}(𝒩)`$ is irreducible for any simple $`g`$, and the dense $`G`$-orbit in it is Richardson.
## 11. On invariants and null-cones for generalised Takiff Lie algebras
Following , we recall the definition of a generalised Takiff Lie algebra. The infinite-dimensional 𝕜-vector space $`q_{\mathrm{}}:=q\text{𝕜}[𝖳]`$ has a natural structure of a Lie algebra such that $`[x𝖳^l,y𝖳^k]=[x,y]𝖳^{l+k}`$. Then $`q_{(n+1)}={\displaystyle \underset{jn+1}{}}q𝖳^j`$ is an ideal of $`q_{\mathrm{}}`$, and the respective quotient is a generalised Takiff Lie algebra, denoted $`qn`$. We also say that $`qn`$ is the $`n`$-th Takiff algebra. Write $`Qn`$ for the corresponding connected group. Clearly, $`dimqn=(n+1)dimq`$ and $`q1qq`$. The main results of are the following:
(i) $`\mathrm{ind}qn=(n+1)\mathrm{ind}q`$,
(ii) if $`q=g`$ is semisimple, then $`\text{𝕜}[gn]^{Gn}`$ is a polynomial algebra whose set of basic invariants is explicitly described.
Actually, the authors of work with invariants of the coadjoint representation of $`Gn`$, but this makes no difference, since $`gn`$ is quadratic.
In this section, we generalise the results from (ii) in the spirit of Section 7. Let $`qn^u`$ denote the image of $`q_1`$ in $`qn`$. It is a nilpotent Lie algebra, which is noncommutative for $`n2`$, and $`qnqqn^u`$. Accordingly, one obtains the semi-direct product structure of the group: $`Qn=QQn^u`$.
11.1 Theorem. Suppose $`q`$ satisfies conditions
(1) $`\text{𝕜}[q]^Q`$ is a polynomial algebra;
(2) $`\mathrm{max}dim_{xq}Qx=dimqdimq//Q`$;
(3) If $`\pi _Q:qq//Q`$ is the quotient morphism and $`\mathrm{\Omega }:=\{xq(\text{d}\pi _Q)_x\text{ is onto }\}`$, then $`q\mathrm{\Omega }`$ contains no divisors.
Then
* $`\text{𝕜}[qn]^{Qn^u}`$ is a polynomial algebra of Krull dimension $`dimq+ndimq//Q`$ whose algebraically independent generators can explicitly be described;
* $`\text{𝕜}[qn]^{Qn}`$ is a polynomial algebra of Krull dimension $`(n+1)dimq//Q`$ whose algebraically independent generators can explicitly be described;
Proof. Let $`𝒙=x_0+ϵx_1+\mathrm{}+ϵ^nx_n`$ denote the image of $`_{i=0}^nx_i𝖳^i`$ in $`qn`$. Here each $`x_iq`$ and $`ϵ`$ is regarded as the image of $`𝖳`$ in $`\text{𝕜}[𝖳]/(𝖳^{n+1})`$ Set $`m=dimq//Q`$, and let $`f_1,\mathrm{},f_m`$ be a set of basic invariants in $`\text{𝕜}[q]^Q`$. Expand the polynomial $`f_i(x_0+ϵx_1+\mathrm{}+ϵ^nx_n)`$ using the relation $`ϵ^{n+1}=0`$. We obtain
$$f_i(x_0+ϵx_1+\mathrm{}+ϵ^nx_n)=\underset{j=0}{\overset{n}{}}ϵ^j\widehat{F}_i^{(j)}(x_0,x_1,\mathrm{},x_n).$$
Following the argument in \[33, Sect. III\], one proves that
$`\widehat{F}_i^{(j)}`$ depends only on $`x_0,\mathrm{},x_j`$ and
(11.2)
$$\widehat{F}_i^{(j)}(x_0,\mathrm{},x_j)=(\text{d}f_i)_{x_0},x_j+p_{ij}(x_0,\mathrm{},x_{j1}).$$
It follows from the construction that all $`\widehat{F}_i^{(j)}`$ belong to $`\text{𝕜}[qn]^{Qn}`$.
(i) Making use of Lemma 6 and Eq. (11.2), we prove that the polynomials $`\widehat{F}_i^{(j)}`$, $`i=1,\mathrm{},m`$, $`j=1,\mathrm{},n`$, and the coordinates on the first factor in $`qn`$ freely generate $`\text{𝕜}[qn]^{Qn^u}`$.
Consider the mapping
$$\psi :qnq\times \text{𝕜}^{nm},$$
given by $`\psi (𝒙)=(x_0,\widehat{F}_1^{(1)}(𝒙),\mathrm{},\widehat{F}_m^{(n)}(𝒙))`$. Here we regard $`q`$ as $`qn/qn^u`$, so that $`q\times \text{𝕜}^{nm}`$ is a variety with trivial $`Qn^u`$-action. If $`x_0\mathrm{\Omega }`$, then $`(\text{d}f_i)_{x_0}`$ are linearly independent. Therefore Eq. (11.2) shows that the system $`\widehat{F}_i^{(j)}(x_0+ϵy_1+\mathrm{}+ϵ^ny_n)=\alpha _i^{(j)}`$, $`i=1,\mathrm{},m`$, $`j=1,\mathrm{},n`$ has a solution, say $`(y_1,\mathrm{},y_n)`$, for any $`nm`$-tuple $`𝜶=(\alpha _1^{(1)},\mathrm{},\alpha _m^{(n)})`$. Indeed, $`(y_1,\mathrm{},y_n)`$ can be computed consecutively: First $`y_1`$, then $`y_2`$, and so on. Hence $`\mathrm{Im}\psi \mathrm{\Omega }\times \text{𝕜}^{nm}`$, i.e., $`\mathrm{Im}\psi `$ contains a big open subset of $`q\times \text{𝕜}^{nm}`$. This also implies that the coordinates on $`q`$ and the polynomials $`\widehat{F}_i^{(j)}`$ are algebraically independent. It follows that
$$\underset{𝒙qn}{\mathrm{max}}dimQn^u𝒙dimqndimqmn=n(dimqm).$$
Next, consider $`\mathrm{\Omega }^{}=\mathrm{\Omega }\{zqdimQz=dimqm\}`$. In view of condition (2), it is still a non-empty open $`Q`$-stable subset of $`q`$. Fix $`x_0\mathrm{\Omega }^{}`$, and let $`(\overline{y}_1,\mathrm{},\overline{y}_n)`$ be a solution to the system $`\widehat{F}_i^{(j)}(x_0+ϵy_1+\mathrm{}+ϵ^ny_n)=\alpha _i^{(j)}`$, $`i=1,\mathrm{},m`$, $`j=1,\mathrm{},n`$. Then $`\psi ^1(x_0,𝜶)Qn^u(x_0+_{i=1}^nϵ^i\overline{y}_i)`$. Since $`x_0\mathrm{\Omega }`$, we have $`dim\psi ^1(x_0,𝜶)=n(dimqm)`$. On the other hand, the following holds
Claim. If $`xq^{reg}`$, then $`dimQn^u(x+ϵy_1+\mathrm{}+ϵ^ny_n)=ndimQx=n(dimqm)`$ for any $`(y_1,\mathrm{},y_n)q^n`$.
Proof of the claim. We argue by induction on $`n`$. For $`n=1`$, the assertion is obvious. Assume that $`n2`$. Consider the $`Qn^u`$-equivariant projection
$$(x+\underset{1}{\overset{n}{}}ϵ^iy_iqn)\stackrel{p}{}(x+\underset{1}{\overset{n1}{}}ϵ^iy_iqn1).$$
Let $`𝒪_n`$ denote the $`Qn^u`$-orbit of $`x+_1^nϵ^iy_i`$. Then $`p(𝒪_n)=𝒪_{n1}`$. By the induction hypothesis, $`dim𝒪_{n1}=(n1)(dimqm)`$. It is easily seen that
$$p^1(x+\underset{1}{\overset{n1}{}}ϵ^iy_i)𝒪_nx+\underset{1}{\overset{n1}{}}ϵ^iy_i+ϵ^n(y_n+[q,x]).$$
For, the right hand side is precisely an orbit of the subgroup $`\mathrm{exp}(ϵ^nq)Qn^u`$. Hence $`dim𝒪_nn(dimqm)`$. But it is already proved that the dimension of every $`Qn^u`$-orbit is at most $`n(dimqm)`$. Hence $`\psi ^1(x_0,𝜶)=Qn^u(x_0+ϵ\overline{y}_1+\mathrm{}+ϵ^n\overline{y}_n)`$ for dimension reason. Thus, a generic fibre of $`\psi `$ is an $`Qn^u`$-orbit, and Lemma 6 applies here.
(ii) Follows from (i) and the description of $`Qn^u`$-invariants. 11.3 Remark. It was noticed in Section 9 that any Takiff algebra $`qq`$ is a $`_2`$-contraction of $`qq`$. Similar phenomenon holds for the generalised Takiff algebras: $`qn`$ is a contraction of $`q\mathrm{}q=(n+1)q`$. The starting point for constructing such a contraction is to consider the action $`_{n+1}`$ on $`(n+1)q`$ that cyclically permutes the summands. On the other hand, given $`qn`$, it can further be contracted to $`qq^n`$, the ”usual” semi-direct product, where $`q^n`$ is regarded as commutative Lie algebra. The details are left to the reader. Thus,
$$q\mathrm{}q=(n+1)qqnqq^n$$
is a chain of contractions. Using Eq. (10.7) and Eq. (11.2) one can easily prove that if $`g`$ is semisimple and $`g2`$ is the second Takiff Lie algebra, then the quotient morphism $`\pi _{G2}:g2g2//G2`$ is equidimensional. This is a particular case of the theorem of Eisenbud-Frenkel mentioned in the introduction. |
warning/0506/quant-ph0506026.html | ar5iv | text | # Standard and Non–Standard Quantum Models : A Non–Commutative Version of the Classical System of SU(2) and SU(1,1) Arising from Quantum Optics
## 1 Introduction
This is a challenging paper including some review of and new results, and our ultimate aim is to construct a unified theory of Non–Commutative (Differential) Geometry and Quantum Computation.
The Hopf bundles (which are famous examples of fiber bundles) over $`𝐊=𝐑`$, $`𝐂`$, $`𝐇`$ (the field of quaternion numbers), $`𝐎`$ (the field of octanion numbers) are classical objects and they are never written down in a local manner. If we write them locally then we are forced to encounter singular lines called the Dirac strings, see , .
It is very interesting to comment that the Hopf bundles correspond to topological solitons called Kink, Monopole, Instanton, Generalized Instanton respectively, see for example , , . Therefore they are very important objects to study.
Berry has given another expression to the Hopf bundle and Dirac strings by making use of a Hamiltonian (a simple spin model including the parameters $`x`$, $`y`$ and $`z`$), see the paper(s) in . We call this the Berry model for simplicity. In this paper let us restrict to the case of $`𝐊`$=$`𝐂`$.
We also construct a pseudo Berry model by replacing the Pauli matrices (the generators of $`su(2)`$) in the Hamiltonian with the generators of $`su(1,1)`$. For this model the “Hamiltonian” is not hermite and a bundle is defined, which is called the pseudo Hopf bundle for simplicity. However, it is topologically trivial (therefore, there are no Dirac strings).
We would like to make the Hopf and pseudo Hopf bundles non–commutative. Whether such a generalization is meaningful or not is not clear at the current time, however it is worth trying, see for example , or more recently and its references.
By the way, we are studying a quantum computation based on Cavity QED and one of the basic tools is the Jaynes–Cummings model (or more generally the Tavis–Cummings one), , , , . This is given as a “half” of the Dicke model under the resonance condition and rotating wave approximation associated to it. If the resonance condition is not taken, then this model gives a non–commutative version of the Berry model. However, this new one is different from usual one because $`x`$ and $`y`$ coordinates are quantized, while $`z`$ coordinate is not.
We also construct a non–commutative version of the pseudo Berry model by replacing the generators as in the classical case. In this case, since the eigenvalues of the pseudo Hamiltonian should be real, the domain is extremely limited in the Fock space.
From the non–commutative Berry model we construct a non–commutative version of the Hopf bundle by making use of so–called Quantum Diagonalization Method developed in . Then we see that the Dirac strings appear in only states containing the ground one ($`\times \{|0\}\{|0\}\times `$), while they don’t appear in excited states ($`\times \times \{|0\}\{|0\}\times `$), where $``$ is the Fock space generated by $`\{a,a^{},N=a^{}a\}`$,
This means that classical singularities are not universal in the process of non–commutativization, which is a very interesting phenomenon. This is one of reasons why we consider non–commutative generalizations (which are not necessarily unique) of classical geometry.
We also construct a non–commutative version of the pseudo Hopf bundle in the non–commutative pseudo Berry model. Since in this case the bundle is trivial and there are no Dirac strings, the situation becomes easy.
Moreover, we construct a non–commutative version of the Veronese mapping which is the mapping from $`𝐂P^1`$ to $`𝐂P^n`$ with mapping degree $`n`$. The mapping degree is usually defined by making use of the (first–) Chern class, so our mapping will become important if a non–commutative (or quantum) “Chern class” would be constructed.
We also construct a non–commutative version of the pseudo Veronese mapping which is the mapping from $`𝐂Q^1`$ to $`𝐂Q^n`$ with mapping degree $`n`$.
We challenge to construct a non–commutative version of the spin representation of group $`SU(2)`$. However, our trial is not enough because we could not construct the general case except for the special cases of spin $`j=1`$ and $`j=3/2`$. In this problem, we meet a difficulty coming from the non–commutativity. Further study constructing a general theory will be required.
We also challenge to construct a non–commutative version of the spin representation of group $`SU(1,1)`$. However, unitary representations are infinite dimensional from the starting point even in the classical case. To develop a unitary theory of non–commutative system of $`SU(1,1)`$ we need an infinite number of non–commutative systems, which means a kind of second non–commutativization. Therefore our trial is not enough, so that further study will be required.
Why do we consider non–commutative versions of classical field models ? What is an advantage to consider such a generalization ? Such natural questions arise. This paper may give one of answers. Moreover, readers will find many interesting (challenging) problems.
For the convenience of readers this paper is arranged as the first subsection is the system based on $`SU(2)`$ and the next one is the system based on $`SU(1,1)`$. This contrast may make the similarity and difference between the standard and non–standard models clear. We also add many appendices to make the text clear.
The contents of the paper are as follows :
Section 1 Introduction
Section 2 Mathematical Preliminaries
2.1 Classical SU(2) System $`\mathrm{}`$ Compact Case
2.2 Classical SU(1,1) System $`\mathrm{}`$ Non-Compact Case
Section 3 Standard and Non-Standard Berry Models and Dirac Strings
3.1 Standard Berry Model and Dirac Strings
3.1 Non-Standard Berry Model
Section 4 Non-Commutative Models Arizing from the Jaynes-Cummings Model
4.1 Standard Quantum Model
4.1 Non-Standard Quantum Model
Section 5 Non-Commutative Hopf and Pseudo Hopf Bundles
5.1 Non-Commutative Hopf Bundle
5.2 Non-Commutative Pseudo Hopf Bundle
Section 6 Non-Commutative Veronese and Pseudo Veronese Mappings
6.1 Non-Commutative Veronese Mapping
6.2 Non-Commutative Pseudo Veronese Mapping
Section 7 Non-Commutative Representation Theory
7.1 Non-Commutative Version of $`SU(2)`$ Case
7.2 Non-Commutative Version of $`SU(1,1)`$ Case
Section 8 Discussion
Appendix
A Classical Theory of Projective Spaces
B Local Coordinate of the Projector
C Some Calculations of First Chern Class
D Difficulty of Tensor Decomposition
E Calculation of Some Integrals
## 2 Mathematical Preliminaries
In this section we prepare some mathematical preliminaries for the following sections.
### 2.1 Classical $`SU(2)`$ System $`\mathrm{}`$ Compact Case
The compact Lie group $`SU(2)`$ and its Lie algebra $`isu(2)`$ ($`i=\sqrt{1}`$) are
$$SU(2)=\left\{AM(2;𝐂)\right|A^{}A=1_2,\text{det}(A)=1\}$$
(1)
and
$$su(2)=\left\{XM(2;𝐂)\right|X^{}=X,\text{tr}(X)=0\}.$$
(2)
The algebra is generated by the famous Pauli matrices $`\sigma _j(j=13)`$
$$\sigma _1=\left(\begin{array}{cc}& 1\\ 1& \end{array}\right),\sigma _2=\left(\begin{array}{cc}& i\\ i& \end{array}\right),\sigma _3=\left(\begin{array}{cc}1& \\ & 1\end{array}\right);1_2=\left(\begin{array}{cc}1& \\ & 1\end{array}\right),$$
and the map $`su(2)SU(2)`$ is given as
$$\underset{j=1}{\overset{3}{}}x_j\sigma _j\text{exp}\left(i\underset{j=1}{\overset{3}{}}x_j\sigma _j\right).$$
We usually use
$$\sigma _+(1/2)(\sigma _1+i\sigma _2)=\left(\begin{array}{cc}0& 1\\ 0& 0\end{array}\right),\sigma _{}(1/2)(\sigma _1i\sigma _2)=\left(\begin{array}{cc}0& 0\\ 1& 0\end{array}\right).$$
Then the $`su(2)`$ relation
$$[\stackrel{~}{\sigma }_3,\sigma _+]=\sigma _+,[\stackrel{~}{\sigma }_3,\sigma _{}]=\sigma _{},[\sigma _+,\sigma _{}]=2\stackrel{~}{\sigma }_3$$
is well–known, where $`\stackrel{~}{\sigma }_3=(1/2)\sigma _3`$.
Let us note that
$$A=\left(\begin{array}{cc}\alpha & \overline{\beta }\\ \beta & \overline{\alpha }\end{array}\right),|\alpha |^2+|\beta |^2=1$$
(3)
is an element in $`SU(2)`$.
### 2.2 Classical $`SU(1,1)`$ System $`\mathrm{}`$ Non–Compact Case
The non–compact Lie group $`SU(1,1)`$ and its Lie algebra $`isu(1,1)`$ are
$$SU(1,1)=\left\{BM(2;𝐂)\right|B^{}JB=J,\text{det}(B)=1\}$$
(4)
and
$$su(1,1)=\left\{YM(2;𝐂)\right|Y^{}=JYJ,\text{tr}(Y)=0\}$$
(5)
where $`J=\sigma _3`$. The algebra is generated by the matrices $`\tau _j(j=13)`$
$$\tau _1=\left(\begin{array}{cc}& 1\\ 1& \end{array}\right),\tau _2=\left(\begin{array}{cc}& i\\ i& \end{array}\right),\tau _3=\left(\begin{array}{cc}1& \\ & 1\end{array}\right)=\sigma _3,$$
and the map $`su(1,1)SU(1,1)`$ is given as
$$\underset{j=1}{\overset{3}{}}x_j\tau _j\text{exp}\left(i\underset{j=1}{\overset{3}{}}x_j\tau _j\right).$$
We usually use
$$\tau _+(1/2)(\tau _1+i\tau _2)=\left(\begin{array}{cc}0& 1\\ 0& 0\end{array}\right),\tau _{}(1/2)(\tau _1i\tau _2)=\left(\begin{array}{cc}0& 0\\ 1& 0\end{array}\right).$$
Then the $`su(1,1)`$ relation
$$[\stackrel{~}{\tau }_3,\tau _+]=\tau _+,[\stackrel{~}{\tau }_3,\tau _{}]=\tau _{},[\tau _+,\tau _{}]=2\stackrel{~}{\tau }_3$$
is well–known, where $`\stackrel{~}{\tau }_3=(1/2)\tau _3`$.
Let us note that
$$B=\left(\begin{array}{cc}\alpha & \overline{\beta }\\ \beta & \overline{\alpha }\end{array}\right),|\alpha |^2|\beta |^2=1$$
(6)
is an element in $`SU(1,1)`$.
## 3 Standard and Non–Standard Berry Models and Dirac Strings
We explain the way which Berry used in to construct the Hopf bundle and Dirac strings corresponding to the compact case, and next construct ones corresponding to the non–compact case.
### 3.1 Standard Berry Model and Dirac Strings
The Hamiltonian used by Berry is a simple spin model
$$H_B=x\sigma _1+y\sigma _2+z\sigma _3=(xiy)\sigma _++(x+iy)\sigma _{}+z\sigma _3=\left(\begin{array}{cc}z& xiy\\ x+iy& z\end{array}\right)$$
(7)
where $`x`$, $`y`$ and $`z`$ are parameters. This Hamiltonian is of course hermite. We would like to diagonalize $`H_B`$ above. The eigenvalues are
$$\lambda =\pm r\pm \sqrt{x^2+y^2+z^2}$$
and corresponding orthonormal eigenvectors are
$$|r=\frac{1}{\sqrt{2r(r+z)}}\left(\begin{array}{c}r+z\\ x+iy\end{array}\right),|r=\frac{1}{\sqrt{2r(r+z)}}\left(\begin{array}{c}x+iy\\ r+z\end{array}\right).$$
Here we assume $`(x,y,z)𝐑^3\{(0,0,0)\}𝐑^3\{0\}`$ to avoid a degenerate case. Therefore a unitary matrix defined by
$$A_I=(|r,|r)=\frac{1}{\sqrt{2r(r+z)}}\left(\begin{array}{cc}r+z& x+iy\\ x+iy& r+z\end{array}\right)$$
(8)
makes $`H_B`$ diagonal like
$$H_B=A_I\left(\begin{array}{cc}r& \\ & r\end{array}\right)A_I^{}A_ID_BA_I^{}.$$
(9)
We note that the unitary matrix $`A_I`$ is not defined on the whole space $`𝐑^3\{0\}`$. The defining region of $`U_I`$ is
$$D_I=𝐑^3\{0\}\{(0,0,z)𝐑^3|z<0\}.$$
(10)
The removed line $`\{(0,0,z)𝐑^3|z<0\}`$ is just the (lower) Dirac string, which is impossible to add to $`D_I`$.
Next, we have another diagonal form of $`H_B`$ like
$$H_B=A_{II}D_BA_{II}^{}$$
(11)
with the unitary matrix $`A_{II}`$ defined by
$$A_{II}=\frac{1}{\sqrt{2r(rz)}}\left(\begin{array}{cc}xiy& r+z\\ rz& x+iy\end{array}\right).$$
(12)
The defining region of $`A_{II}`$ is
$$D_{II}=𝐑^3\{0\}\{(0,0,z)𝐑^3|z>0\}.$$
(13)
The removed line $`\{(0,0,z)𝐑^3|z>0\}`$ is the (upper) Dirac string, which is also impossible to add to $`D_{II}`$.
Here we have diagonalizations of two types for $`H_B`$, so a natural question comes about. What is a relation between $`A_I`$ and $`A_{II}`$ ? If we define
$$\mathrm{\Phi }=\frac{1}{\sqrt{x^2+y^2}}\left(\begin{array}{cc}xiy& \\ & x+iy\end{array}\right)$$
(14)
then it is easy to see
$$A_{II}=A_I\mathrm{\Phi }.$$
We note that $`\mathrm{\Phi }`$ (which is called a transition function) is not defined on the whole $`z`$–axis.
What we would like to emphasize here is that the diagonalization of $`H_B`$ is not given globally (on $`𝐑^3\{0\}`$). However, the dynamics is perfectly controlled by the system
$$\left\{(A_I,D_I),(A_{II},D_{II}),\mathrm{\Phi },D_ID_{II}=𝐑^3\{0\}\right\},$$
(15)
which defines a famous fiber bundle called the Hopf bundle associated to the complex numbers $`𝐂`$ <sup>1</sup><sup>1</sup>1The base space $`𝐑^3\{0\}`$ is homotopic to the two–dimensional sphere $`S^2`$,
$$S^1S^3S^2,$$
see .
The projector corresponding to the Hopf bundle is given as
$$P(x,y,z)=A_IP_0A_I^{}=A_{II}P_0A_{II}^{}=\frac{1}{2r}\left(\begin{array}{cc}r+z& xiy\\ x+iy& rz\end{array}\right),$$
(16)
where $`P_0`$ is the basic one
$$P_0=\left(\begin{array}{cc}1& \\ & 0\end{array}\right)M(2,𝐂).$$
It is well–known that $`P`$ satisfies the relations
$$1)P^2=P,2)P=P^{},3)\text{tr}P=1.$$
We note that in (16) Dirac strings don’t appear because the projector $`P`$ is expressed globally.
### 3.2 Non–Standard Berry Model
The “Hamiltonian” that we consider here is a modified one of the Berry model
$$H_{pB}=x\tau _1+y\tau _2+z\tau _3=(xiy)\tau _++(x+iy)\tau _{}+z\tau _3=\left(\begin{array}{cc}z& xiy\\ (x+iy)& z\end{array}\right)$$
(17)
where $`x`$, $`y`$ and $`z`$ are parameters. This is not hermite. As a tentative terminology we call this a pseudo Berry model. We would like to diagonalize $`H_{pB}`$. The eigenvalues are
$$\lambda =\pm s\pm \sqrt{z^2x^2y^2},$$
so the defining domain is
$$D\left\{(x,y,z)𝐑^3\right|z^2x^2y^2>0\}$$
Here, to avoid a degenerate case of eigenvalues we removed the case of $`z^2x^2y^2=0`$. We note that $`D`$ is not connected and consists of two domains $`D_+`$ and $`D_{}`$ defined by
$$D_+=\left\{(x,y,z)D\right|z>0\}\text{and}D_{}=\left\{(x,y,z)D\right|z<0\}.$$
(18)
The corresponding orthonormal eigenvectors are
$$|s=\frac{1}{\sqrt{2s(s+z)}}\left(\begin{array}{c}r+z\\ (x+iy)\end{array}\right),|s=\frac{1}{\sqrt{2s(s+z)}}\left(\begin{array}{c}x+iy\\ s+z\end{array}\right).$$
Therefore a matrix defined by
$$B_I=(|s,|s)=\frac{1}{\sqrt{2s(s+z)}}\left(\begin{array}{cc}s+z& x+iy\\ (x+iy)& s+z\end{array}\right)$$
(19)
makes $`H_{pB}`$ diagonal like
$$H_{pB}=B_I\left(\begin{array}{cc}s& \\ & s\end{array}\right)B_I^1B_ID_{pB}B_I^1.$$
(20)
We note that the matrix $`B_I`$ is an element of the non–compact group $`SU(1,1)`$, and is not defined on $`D_{}`$ because $`s+z<0`$. Moreover, $`B_I`$ is defined on the whole $`D_I=D_+`$, so there is no singular line like Dirac strings.
A comment is in order. We have another diagonal form of $`H_{pB}`$ like
$$H_{pB}=B_{II}D_{pB}B_{II}^1$$
(21)
with the matrix $`B_{II}`$ in $`SU(1,1)`$ defined by
$$B_{II}=\frac{1}{\sqrt{2s(zs)}}\left(\begin{array}{cc}x+iy& zs\\ zs& (x+iy)\end{array}\right).$$
(22)
The defining region of $`B_{II}`$ is
$$D_{II}=D_+\{(0,0,z)D_+\}.$$
(23)
The removed line $`\{(0,0,z)D_+\}`$ is the (upper) Dirac string, which is also impossible to add to $`D_{II}`$.
However, with a singular transformation $`\mathrm{\Phi }`$ (not defined on the whole $`z`$–axis) defined by
$$\mathrm{\Phi }=\frac{1}{\sqrt{x^2+y^2}}\left(\begin{array}{cc}xiy& \\ & x+iy\end{array}\right)$$
(24)
we can remove it because
$$B_I=B_{II}\mathrm{\Phi }.$$
What we would like to emphasize in this case is that the diagonalization of $`H_{pB}`$ is given globally on $`D_I`$, which is very different from the compact case.
$$\{H_{pB},B_I,D_I\}.$$
(25)
Here, as a tentative terminology we call this system a pseudo Hopf bundle corresponding to the Hopf bundle in the compact case. However, this doesn’t define a topological object because the domain $`D_I`$ is contractible (trivial in the sense of topology).
The projector corresponding to the case is given as
$$Q(x,y,z)=B_IQ_0B_I^1=\frac{1}{2s}\left(\begin{array}{cc}z+s& xiy\\ (x+iy)& z+s\end{array}\right),$$
(26)
where $`Q_0=P_0`$.
$`Q`$ satisfies the relations
$$1)Q^2=Q,2)JQJ=Q^{},3)\text{tr}Q=1.$$
In the following we omit the suffix $`I`$ for simplicity.
## 4 Non–Commutative Models Arising from the <br>Jaynes–Cummings Model
In this section let us explain the Jaynes–Cummings model which is well–known in quantum optics, see , . From this we obtain a standard quantum model which is a natural extension of the (classical) Berry model. On the other hand, we obtain a non–standard quantum model by replacing the bases of $`su(2)`$ with the bases of $`su(1,1)`$.
### 4.1 Standard Quantum Model
The Hamiltonian of Jaynes–Cummings model can be written as follows (we set $`\mathrm{}=1`$ for simplicity)
$$H=\omega 1_2a^{}a+\frac{\mathrm{\Delta }}{2}\sigma _3\mathrm{𝟏}+g\left(\sigma _+a+\sigma _{}a^{}\right),$$
(27)
where $`\omega `$ is the frequency of single radiation field, $`\mathrm{\Delta }`$ the energy difference of two level atom, $`a`$ and $`a^{}`$ are annihilation and creation operators of the field, and $`g`$ a coupling constant. We assume that $`g`$ is small enough (a weak coupling regime). See the figure 3 as an image of the Jaynes–Cummings model (we don’t repeat here).
Now we consider the evolution operator of the model. We rewrite the Hamiltonian (27) as follows.
$$H=\omega 1_2a^{}a+\frac{\omega }{2}\sigma _3\mathrm{𝟏}+\frac{\mathrm{\Delta }\omega }{2}\sigma _3\mathrm{𝟏}+g\left(\sigma _+a+\sigma _{}a^{}\right)H_1+H_2.$$
(28)
Then it is easy to see $`[H_1,H_2]=0`$, which leads to $`\text{e}^{itH}=\text{e}^{itH_1}\text{e}^{itH_2}.`$
In the following we consider $`\text{e}^{itH_2}`$ in which the resonance condition $`\mathrm{\Delta }\omega =0`$ is not taken. For simplicity we set $`\theta =\frac{\mathrm{\Delta }\omega }{2g}(0)`$ <sup>2</sup><sup>2</sup>2Since the Jaynes–Cummings model is obtained by the Dicke model under some resonance condition on parameters included, it is nothing but an approximate one in the neighborhood of the point, so we must assume that $`|\theta |`$ is small enough. However, as a model in mathematical physics there is no problem to take $`\theta `$ be arbitrary then
$$H_2=g\left(\sigma _+a+\sigma _{}a^{}+\frac{\mathrm{\Delta }\omega }{2g}\sigma _3\mathrm{𝟏}\right)=g\left(\sigma _+a+\sigma _{}a^{}+\theta \sigma _3\mathrm{𝟏}\right).$$
For further simplicity we set
$$H_{JC}=\sigma _+a+\sigma _{}a^{}+\theta \sigma _3\mathrm{𝟏}=\left(\begin{array}{cc}\theta & a\\ a^{}& \theta \end{array}\right),[a,a^{}]=\mathrm{𝟏}$$
(29)
where we have written $`\theta `$ in place of $`\theta \mathrm{𝟏}`$ for simplicity.
$`H_{JC}`$ can be considered as a non-commutative version of $`H_B`$ under the correspondence $`axiy,a^{}x+iy\text{and}\theta z`$ :
$$H_B=\left(\begin{array}{cc}z& xiy\\ x+iy& z\end{array}\right),[xiy,x+iy]=0H_{JC}=\left(\begin{array}{cc}\theta & a\\ a^{}& \theta \end{array}\right),[a,a^{}]=\mathrm{𝟏}.$$
(30)
That is, $`x`$ and $`y`$ coordinates are quantized, while $`z`$ coordinate is not, which is different from usual one, see for example . It may be possible for us to call this a non–commutative Berry model. We note that this model is derived not “by hand” but by the model in quantum optics itself.
### 4.2 Non–Standard Quantum Model
Similarly, from (29) we can define
$$H_{pJC}=\tau _+a+\tau _{}a^{}+\theta \tau _3\mathrm{𝟏}=\left(\begin{array}{cc}\theta & a\\ a^{}& \theta \end{array}\right),[a,a^{}]=\mathrm{𝟏}$$
(31)
by replacing $`\{\sigma _j\}`$ with $`\{\tau _j\}`$. In this case this model is derived “by hand”. It satisfies the $`su(1,1)`$ like relation (see (5))
$$𝐉H_{pJC}𝐉=H_{pJC}^{};𝐉=\left(\begin{array}{cc}\mathrm{𝟏}& \\ & \mathrm{𝟏}\end{array}\right)=J\mathrm{𝟏}.$$
$`H_{pJC}`$ can be considered as a non-commutative version of $`H_{pB}`$ under the correspondence $`axiy,a^{}x+iy\text{and}\theta z`$ :
$$H_{pB}=\left(\begin{array}{cc}z& xiy\\ (x+iy)& z\end{array}\right),[xiy,x+iy]=0H_{pJC}=\left(\begin{array}{cc}\theta & a\\ a^{}& \theta \end{array}\right),[a,a^{}]=\mathrm{𝟏}.$$
(32)
A comment is in order. In place of the Hamiltonian (27) we can consider the following pseudo Hamiltonian
$$H_p=\omega 1_2a^{}a+\frac{\mathrm{\Delta }}{2}\tau _3\mathrm{𝟏}+g\left(\tau _+a+\tau _{}a^{}\right)$$
(33)
by replacing $`\{\sigma _+,\sigma _{},\sigma _3\}`$ with $`\{\tau _+,\tau _{},\tau _3\}`$. This is not hermite (namely, not a convensional one), so we don’t know whether this model is useful or not at the current time. It is interesting to note that the model has been considered by .
In a forthcoming paper we will extend this “Hamiltonian” and determine its structure in detail like .
## 5 Non–Commutative Hopf and Pseudo Hopf Bundles
In this section we construct a non–commutative version of the Hopf and pseudo Hopf bundles by making (29) and (31) diagonal, which is a “natural” extension in the section 3.
First of all let us recall a Fock space. For $`a`$ and $`a^{}`$ we set $`Na^{}a`$ which is called the number operator, then we have
$$[N,a^{}]=a^{},[N,a]=a,[a^{},a]=\mathrm{𝟏}.$$
(34)
Let $``$ be the Fock space generated by $`\{a,a^{},N\}`$
$$=\text{Vect}_𝐂\{|0,|1,\mathrm{},|n,\mathrm{}\}.$$
(35)
The actions of $`a`$ and $`a^{}`$ on $``$ are given by
$$a|n=\sqrt{n}|n1,a^{}|n=\sqrt{n+1}|n+1,N|n=n|n$$
(36)
where $`|0`$ is a normalized vacuum ($`a|0=0\mathrm{and}0|0=1`$). From (36) state $`|n`$ for $`n1`$ are given by
$$|n=\frac{(a^{})^n}{\sqrt{n!}}|0.$$
(37)
These states satisfy the orthogonality and completeness conditions
$$m|n=\delta _{mn}\text{and}\underset{n=0}{\overset{\mathrm{}}{}}|nn|=\mathrm{𝟏}.$$
(38)
### 5.1 Non–Commutative Hopf Bundle
First we make the Hamiltonian (29) diagonal like in Section 2 and research whether “Dirac strings” exist or not in this non–commutative model, which is very interesting from not only quantum optical but also mathematical point of view.
It is easy to see
$$H_{JC}=\left(\begin{array}{cc}\theta & a\\ a^{}& \theta \end{array}\right)=\left(\begin{array}{cc}1& \\ & a^{}\frac{1}{\sqrt{N+1}}\end{array}\right)\left(\begin{array}{cc}\theta & \sqrt{N+1}\\ \sqrt{N+1}& \theta \end{array}\right)\left(\begin{array}{cc}1& \\ & \frac{1}{\sqrt{N+1}}a\end{array}\right)$$
(39)
from . Then the middle matrix in the right hand side can be considered as a classical one, so we can diagonalize it easily
$$\left(\begin{array}{cc}\theta & \sqrt{N+1}\\ \sqrt{N+1}& \theta \end{array}\right)=\{\begin{array}{c}A_I\left(\begin{array}{cc}R(N+1)& \\ & R(N+1)\end{array}\right)A_I^{}\hfill \\ A_{II}\left(\begin{array}{cc}R(N+1)& \\ & R(N+1)\end{array}\right)A_{II}^{}\hfill \end{array}$$
(40)
where
$$R(N)=\sqrt{N+\theta ^2}$$
and $`A_I`$, $`A_{II}`$ are defined by
$`A_I`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2R(N+1)(R(N+1)+\theta )}}}\left(\begin{array}{cc}R(N+1)+\theta & \sqrt{N+1}\\ \sqrt{N+1}& R(N+1)+\theta \end{array}\right),`$ (43)
$`A_{II}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2R(N+1)(R(N+1)\theta )}}}\left(\begin{array}{cc}\sqrt{N+1}& R(N+1)+\theta \\ R(N+1)\theta & \sqrt{N+1}\end{array}\right).`$ (46)
Now let us rewrite (39) by making use of (40) with (43). Inserting the identity
$$\left(\begin{array}{cc}1& \\ & \frac{1}{\sqrt{N+1}}a\end{array}\right)\left(\begin{array}{cc}1& \\ & a^{}\frac{1}{\sqrt{N+1}}\end{array}\right)=\left(\begin{array}{cc}1& \\ & 1\end{array}\right)$$
gives
$`H_{JC}=`$ $`\left(\begin{array}{cc}1& \\ & a^{}\frac{1}{\sqrt{N+1}}\end{array}\right)A_I\left(\begin{array}{cc}R(N+1)& \\ & R(N+1)\end{array}\right)A_I^{}\left(\begin{array}{cc}1& \\ & \frac{1}{\sqrt{N+1}}a\end{array}\right)`$ (53)
$`=`$ $`\left(\begin{array}{cc}1& \\ & a^{}\frac{1}{\sqrt{N+1}}\end{array}\right)A_I\left(\begin{array}{cc}1& \\ & \frac{1}{\sqrt{N+1}}a\end{array}\right)\left(\begin{array}{cc}1& \\ & a^{}\frac{1}{\sqrt{N+1}}\end{array}\right)\left(\begin{array}{cc}R(N+1)& \\ & R(N+1)\end{array}\right)\times `$ (69)
$`\left(\begin{array}{cc}1& \\ & \frac{1}{\sqrt{N+1}}a\end{array}\right)\left(\begin{array}{cc}1& \\ & a^{}\frac{1}{\sqrt{N+1}}\end{array}\right)A_I^{}\left(\begin{array}{cc}1& \\ & \frac{1}{\sqrt{N+1}}a\end{array}\right)`$
$`=`$ $`U_I\left(\begin{array}{cc}R(N+1)& \\ & R(N)\end{array}\right)U_I^{},`$ (72)
where
$`U_I`$ $`=`$ $`\left(\begin{array}{cc}\frac{1}{\sqrt{2R(N+1)(R(N+1)+\theta )}}& \\ & \frac{1}{\sqrt{2R(N)(R(N)+\theta )}}\end{array}\right)\left(\begin{array}{cc}R(N+1)+\theta & a\\ a^{}& R(N)+\theta \end{array}\right)`$ (77)
$`=`$ $`\left(\begin{array}{cc}R(N+1)+\theta & a\\ a^{}& R(N)+\theta \end{array}\right)\left(\begin{array}{cc}\frac{1}{\sqrt{2R(N+1)(R(N+1)+\theta )}}& \\ & \frac{1}{\sqrt{2R(N)(R(N)+\theta )}}\end{array}\right).`$ (82)
Similarly, we can rewrite (39) by making use of (40) with (46). By inserting the identity
$$\left(\begin{array}{cc}\frac{1}{\sqrt{N+1}}a& \\ & 1\end{array}\right)\left(\begin{array}{cc}a^{}\frac{1}{\sqrt{N+1}}& \\ & 1\end{array}\right)=\left(\begin{array}{cc}1& \\ & 1\end{array}\right)$$
we obtain
$$H_{JC}=U_{II}\left(\begin{array}{cc}R(N)& \\ & R(N+1)\end{array}\right)U_{II}^{},$$
(83)
where
$`U_{II}`$ $`=`$ $`\left(\begin{array}{cc}\frac{1}{\sqrt{2R(N+1)(R(N+1)\theta )}}& \\ & \frac{1}{\sqrt{2R(N)(R(N)\theta )}}\end{array}\right)\left(\begin{array}{cc}a& R(N+1)+\theta \\ R(N)\theta & a^{}\end{array}\right)`$ (88)
$`=`$ $`\left(\begin{array}{cc}a& R(N+1)+\theta \\ R(N)\theta & a^{}\end{array}\right)\left(\begin{array}{cc}\frac{1}{\sqrt{2R(N)(R(N)\theta )}}& \\ & \frac{1}{\sqrt{2R(N+1)(R(N+1)\theta )}}\end{array}\right).`$ (93)
Tidying up these we have
$$H_{JC}=\{\begin{array}{c}U_I\left(\begin{array}{cc}R(N+1)& \\ & R(N)\end{array}\right)U_I^{}\hfill \\ U_{II}\left(\begin{array}{cc}R(N)& \\ & R(N+1)\end{array}\right)U_{II}^{}\hfill \end{array}$$
(94)
with $`U_I`$ and $`U_{II}`$ above. From the equations
$$R(N+1)|0=\sqrt{1+\theta ^2}>\theta ,R(N)|0=\sqrt{\theta ^2}=|\theta |$$
we know
$$\left(R(N)\pm \theta \right)|0=\left(|\theta |\pm \theta \right)|0,$$
so the strings corresponding to Dirac ones exist in only states $`\times \{|0\}\{|0\}\times `$ where $``$ is the Fock space, while in other excited states $`\times \times \{|0\}\{|0\}\times `$ they don’t exist <sup>3</sup><sup>3</sup>3We have identified $`\times `$ with the space of $`2`$–component vectors over $``$, see the figure 3. The phenomenon is very interesting. For simplicity we again call these strings Dirac ones in the following.
The “parameter space” of $`H_{JC}`$ can be identified with $`\times \times 𝐑(,,\theta )`$, so the domains $`D_I`$ of $`U_I`$ and $`D_{II}`$ of $`U_{II}`$ are respectively
$`D_I`$ $`=`$ $`\times \times 𝐑\times \{|0\}\times 𝐑_0,`$ (95)
$`D_{II}`$ $`=`$ $`\times \times 𝐑\left(\times \{|0\}\{|0\}\times \right)\times 𝐑_0`$ (96)
by (77) and (88). We note that
$$D_ID_{II}=\times \times 𝐑\times \{|0\}\times \{\theta =0\}.$$
Then the transition “function” (operator) is given by
$$\mathrm{\Phi }_{JC}=\left(\begin{array}{cc}a\frac{1}{\sqrt{N}}& \\ & \frac{1}{\sqrt{N}}a^{}\end{array}\right)=\left(\begin{array}{cc}\frac{1}{\sqrt{N+1}}a& \\ & a^{}\frac{1}{\sqrt{N+1}}\end{array}\right).$$
Therefore the system
$$\{(U_I,D_I),(U_{II},D_{II}),\mathrm{\Phi }_{JC},D_ID_{II}\}$$
(97)
is a non-commutative version of the Hopf bundle (15). The projector in this case becomes
$`P_{JC}`$ $`=`$ $`U_I\left(\begin{array}{cc}\mathrm{𝟏}& \\ & \mathrm{𝟎}\end{array}\right)U_I^{}=U_{II}\left(\begin{array}{cc}\mathrm{𝟏}& \\ & \mathrm{𝟎}\end{array}\right)U_{II}^{}`$ (102)
$`=`$ $`\{\begin{array}{c}\left(\begin{array}{cc}\frac{1}{2R(N+1)}& \\ & \frac{1}{2R(N)}\end{array}\right)\left(\begin{array}{cc}R(N+1)+\theta & a\\ a^{}& R(N)\theta \end{array}\right)\hfill \\ \left(\begin{array}{cc}R(N+1)+\theta & a\\ a^{}& R(N)\theta \end{array}\right)\left(\begin{array}{cc}\frac{1}{2R(N+1)}& \\ & \frac{1}{2R(N)}\end{array}\right).\hfill \end{array}`$ (113)
Note that the projector $`P_{JC}`$ is not defined on $`\times \{|0\}\times \{\theta =0\}`$ = $`\times \times 𝐑D_ID_{II}`$.
A comment is in order. From (102) we obtain a quantum version of (classical) spectral decomposition (a “quantum spectral decomposition” by Suzuki )
$$H_{JC}=\left(\begin{array}{cc}R(N+1)& \\ & R(N)\end{array}\right)P_{JC}\left(\begin{array}{cc}R(N+1)& \\ & R(N)\end{array}\right)(\mathrm{𝟏}_2P_{JC}).$$
(114)
As a bonus of the decomposition let us rederive the calculation of $`\text{e}^{igtH_{JC}}`$ which has been given in . The result is
$$\text{e}^{igtH_{JC}}=\left(\begin{array}{cc}\text{cos}(tgR(N+1))i\theta \frac{\text{sin}(tgR(N+1))}{R(N+1)}& i\frac{\text{sin}(tgR(N+1))}{R(N+1)}a\\ i\frac{\text{sin}(tgR(N))}{R(N)}a^{}& \text{cos}(tgR(N))+i\theta \frac{\text{sin}(tgR(N))}{R(N)}\end{array}\right)$$
(115)
by making use of (94) (or (114)). We leave it to the readers.
### 5.2 Non–Commutative Pseudo Hopf Bundle
Similarly, we make the Hamiltonian (31) diagonal like in the preceding subsection to study what a non–commutative version of Dirac strings is.
It is easy to see
$$H_{pJC}=\left(\begin{array}{cc}\theta & a\\ a^{}& \theta \end{array}\right)=\left(\begin{array}{cc}1& \\ & a^{}\frac{1}{\sqrt{N+1}}\end{array}\right)\left(\begin{array}{cc}\theta & \sqrt{N+1}\\ \sqrt{N+1}& \theta \end{array}\right)\left(\begin{array}{cc}1& \\ & \frac{1}{\sqrt{N+1}}a\end{array}\right).$$
(116)
The middle matrix in the right hand side can be considered as a classical one, so we can diagonalize it easily
$$\left(\begin{array}{cc}\theta & \sqrt{N+1}\\ \sqrt{N+1}& \theta \end{array}\right)=B\left(\begin{array}{cc}S(N+1)& \\ & S(N+1)\end{array}\right)B^1$$
(117)
where
$$S(N)=\sqrt{\theta ^2N}$$
and $`B`$ defined by
$$B=\frac{1}{\sqrt{2S(N+1)(\theta +S(N+1))}}\left(\begin{array}{cc}S(N+1)+\theta & \sqrt{N+1}\\ \sqrt{N+1}& S(N+1)+\theta \end{array}\right).$$
(118)
In this case the situation changes in a drastic manner. Since $`S(N+1)=\sqrt{\theta ^2(N+1)}`$ where $`N=a^{}a`$ is the number operator, it is clear that only a restricted subspace of the Fock space $``$
$$_n=\text{Vect}_𝐂\{|0,|1,\mathrm{},|n1\}$$
is available if $`n<\theta ^2n+1`$. Moreover, if $`0<\theta ^21`$ there is no subspace that $`S(N+1)`$ is defined !
Similarly in the preceding subsection we have
$$H_{pJC}=V\left(\begin{array}{cc}S(N+1)& \\ & S(N)\end{array}\right)V^1$$
(119)
with $`V`$ defined by
$`V`$ $`=`$ $`\left(\begin{array}{cc}\frac{1}{\sqrt{2S(N+1)(S(N+1)+\theta )}}& \\ & \frac{1}{\sqrt{2S(N)(S(N)+\theta )}}\end{array}\right)\left(\begin{array}{cc}S(N+1)+\theta & a\\ a^{}& S(N)+\theta \end{array}\right)`$ (124)
$`=`$ $`\left(\begin{array}{cc}R(N+1)+\theta & a\\ a^{}& R(N)+\theta \end{array}\right)\left(\begin{array}{cc}\frac{1}{\sqrt{2S(N+1)(S(N+1)+\theta )}}& \\ & \frac{1}{\sqrt{2S(N)(S(N)+\theta )}}\end{array}\right).`$ (129)
We note that $`V`$ above satisfies the relation
$$V^{}𝐉V=𝐉,\text{where}𝐉=\left(\begin{array}{cc}\mathrm{𝟏}& \\ & \mathrm{𝟏}\end{array}\right)V^1=𝐉V^{}𝐉.$$
The “parameter space” of $`H_{pJC}`$ can be identified with
$$\underset{n𝐍}{}_n\times _{n+1}\times \{\theta 𝐑_{>0}|n<\theta ^2n+1\}(,,\theta ).$$
The projector in this case becomes
$`Q_{pJC}`$ $`=`$ $`V\left(\begin{array}{cc}\mathrm{𝟏}& \\ & \mathrm{𝟎}\end{array}\right)V^1`$ (132)
$`=`$ $`\{\begin{array}{c}\left(\begin{array}{cc}\frac{1}{2S(N+1)}& \\ & \frac{1}{2S(N)}\end{array}\right)\left(\begin{array}{cc}S(N+1)+\theta & a\\ a^{}& S(N)\theta \end{array}\right)\hfill \\ \left(\begin{array}{cc}S(N+1)+\theta & a\\ a^{}& S(N)\theta \end{array}\right)\left(\begin{array}{cc}\frac{1}{2S(N+1)}& \\ & \frac{1}{2S(N)}\end{array}\right).\hfill \end{array}`$ (143)
It is easy to see the relations
$$Q_{pJC}^2=Q_{pJC},𝐉Q_{pJC}𝐉=Q_{pJC}^{}.$$
A comment is in order. From (132) we obtain a quantum version of (classical) spectral decomposition
$$H_{pJC}=\left(\begin{array}{cc}S(N+1)& \\ & S(N)\end{array}\right)Q_{pJC}\left(\begin{array}{cc}S(N+1)& \\ & S(N)\end{array}\right)(\mathrm{𝟏}_2Q_{pJC}).$$
(144)
As a bonus of the decomposition let us rederive the calculation of $`\text{e}^{igtH_{pJC}}`$ which seems to be new. The result is
$$\text{e}^{igtH_{pJC}}=\left(\begin{array}{cc}\text{cos}(tgS(N+1))i\theta \frac{\text{sin}(tgS(N+1))}{S(N+1)}& i\frac{\text{sin}(tgS(N+1))}{S(N+1)}a\\ i\frac{\text{sin}(tgS(N))}{S(N)}a^{}& \text{cos}(tgS(N))+i\theta \frac{\text{sin}(tgS(N))}{S(N)}\end{array}\right)$$
(145)
by making use of (119) (or (144)). We leave it to the readers.
We note once more that $`\text{e}^{igtH_{pJC}}`$ is not unitary, but satisfies the relation
$$\left(\text{e}^{igtH_{pJC}}\right)^{}𝐉\text{e}^{igtH_{pJC}}=𝐉.$$
## 6 Non–Commutative Veronese and Pseudo Veronese <br>Mappings
In this section we construct a non–commutative version of the (classical) Veronese Mapping and its noncompact counterpart.
### 6.1 Non–Commutative Veronese Mapping
Let us make a brief review of the Veronese mapping. The map
$$𝐂P^1𝐂P^n$$
is defined as
$$[z_1:z_2][z_1^n:\sqrt{{}_{n}{}^{}C_{1}^{}}z_1^{n1}z_2:\mathrm{}:\sqrt{{}_{n}{}^{}C_{j}^{}}z_1^{nj}z_2^j:\mathrm{}:\sqrt{{}_{n}{}^{}C_{n1}^{}}z_1z_2^{n1}:z_2^n]$$
by making use of the homogeneous coordinate, see Appendix A. We also have another expression of this map by using
$$S_𝐂^1S_𝐂^n:v_1\left(\begin{array}{c}z_1\\ z_2\end{array}\right)v_n\left(\begin{array}{c}z_1^n\\ \sqrt{{}_{n}{}^{}C_{1}^{}}z_1^{n1}z_2\\ \mathrm{}\\ \sqrt{{}_{n}{}^{}C_{j}^{}}z_1^{nj}z_2^j\\ \mathrm{}\\ \sqrt{{}_{n}{}^{}C_{n1}^{}}z_1z_2^{n1}\\ z_2^n\end{array}\right),|z_1|^2+|z_2|^2=1$$
where $`S_𝐂^m=\left\{(w_1,w_2,\mathrm{},w_{m+1})^\text{T}𝐂^{m+1}\right|_{j=1}^{m+1}|w_j|^2=1\}S^{2m+1}`$ and $`𝐂P^m=S_𝐂^m/U(1)`$. Then the Veronese mapping is also written as
$$𝐂P^1𝐂P^n:P_1=v_1v_1^{}P_n=v_nv_n^{}.$$
by using projectors, which is easy to understand.
Moreover, the local map ($`zz_2/z_1`$) is given as
$$𝐂𝐂^n:z\left(\begin{array}{c}\sqrt{{}_{n}{}^{}C_{1}^{}}z\\ \mathrm{}\\ \sqrt{{}_{n}{}^{}C_{j}^{}}z^j\\ \mathrm{}\\ \sqrt{{}_{n}{}^{}C_{n1}^{}}z^{n1}\\ z^n\end{array}\right).$$
See the following picture as a whole.
Next we want to consider a non–commutative version of the map. If we set
$$𝒜\left(\begin{array}{c}X_0\\ Y_0\end{array}\right)=\left(\begin{array}{c}\frac{R(N+1)+\theta }{\sqrt{2R(N+1)(R(N+1)+\theta )}}\\ \frac{1}{\sqrt{2R(N)(R(N)+\theta )}}a^{}\end{array}\right)$$
(146)
from $`U_I`$ in (77), then
$$𝒜^{}𝒜=X_0^2+Y_0^{}Y_0=\mathrm{𝟏}\text{and}Y_0X_0^1=\frac{1}{R(N)+\theta }a^{}Z.$$
That is, $`𝒜=(X_0,Y_0)^T`$ is a non–commutative sphere and $`Z`$ is a kind of “stereographic projection” of the sphere. It is easy to see the following
$$\mathrm{𝟏}+Z^{}Z=\frac{2R(N+1)}{R(N+1)+\theta }=X_0^2X_0=\left(\mathrm{𝟏}+Z^{}Z\right)^{1/2}.$$
(147)
Here let us introduce new notations for the following. For $`j0`$ we set
$`X_j`$ $`=`$ $`{\displaystyle \frac{R(N+1j)+\theta }{\sqrt{2R(N+1j)(R(N+1j)+\theta )}}},`$ (148)
$`Y_j`$ $`=`$ $`\sqrt{{\displaystyle \frac{Nj}{N}}}{\displaystyle \frac{1}{\sqrt{2R(Nj)(R(Nj)+\theta )}}}a^{}.`$ (149)
We list some useful formulas.
$$X_j^2+Y_j^{}Y_j=\mathrm{𝟏}\text{and}Y_j^{}Y_j=Y_{j+1}Y_{j+1}^{}\text{for}j0.$$
(150)
Now we are in a position to define a quantum version of the Veronese mapping which plays a very important role in “classical” Mathematics.
$$𝒜=\left(\begin{array}{c}X_0\\ Y_0\end{array}\right)𝒜_n=\left(\begin{array}{c}X_0^n\\ \sqrt{{}_{n}{}^{}C_{1}^{}}Y_0X_0^{n1}\\ \mathrm{}\\ \sqrt{{}_{n}{}^{}C_{j}^{}}Y_{(j1)}Y_{(j2)}\mathrm{}Y_1Y_0X_0^{nj}\\ \mathrm{}\\ \sqrt{{}_{n}{}^{}C_{n1}^{}}Y_{(n2)}Y_{(n3)}\mathrm{}Y_1Y_0X_0\\ Y_{(n1)}Y_{(n2)}Y_{(n3)}\mathrm{}Y_1Y_0\end{array}\right).$$
(151)
Then it is not difficult to see
$$𝒜_n^{}𝒜_n=\left(X_0^2+Y_0^{}Y_0\right)^n=\left(𝒜^{}𝒜\right)^n=\mathrm{𝟏}.$$
From this we can define the projectors which correspond to projective spaces like
$$𝒫_n=𝒜_n𝒜_n^{},𝒫_1=𝒜𝒜^{},$$
(152)
so the map
$$𝒫_1𝒫_n$$
(153)
is a non-commutative version of the Veronese mapping.
Next, we define a local “coordinate” of the Veronese mapping defined above.
$`𝒜_n`$ $`=`$ $`\left(\begin{array}{c}\mathrm{𝟏}\\ \sqrt{{}_{n}{}^{}C_{1}^{}}Y_0X_0^1\\ \mathrm{}\\ \sqrt{{}_{n}{}^{}C_{j}^{}}Y_{(j1)}Y_{(j2)}\mathrm{}Y_1Y_0X_0^j\\ \mathrm{}\\ \sqrt{{}_{n}{}^{}C_{n1}^{}}Y_{(n2)}Y_{(n3)}\mathrm{}Y_1Y_0X_0^{(n1)}\\ Y_{(n1)}Y_{(n2)}Y_{(n3)}\mathrm{}Y_1Y_0X_0^n\end{array}\right)X_0^n`$
$`=`$ $`\mathrm{}`$
$`=`$ $`\left(\begin{array}{c}\mathrm{𝟏}\\ \sqrt{{}_{n}{}^{}C_{1}^{}}Y_0X_0^1\\ \mathrm{}\\ \sqrt{{}_{n}{}^{}C_{j}^{}}Y_{(j1)}X_{(j1)}^1Y_{(j2)}X_{(j2)}^1\mathrm{}Y_1X_1^1Y_0X_0^1\\ \mathrm{}\\ \sqrt{{}_{n}{}^{}C_{n1}^{}}Y_{(n2)}X_{(n2)}^1Y_{(n3)}X_{(n3)}^1\mathrm{}Y_1X_1^1Y_0X_0^1\\ Y_{(n1)}X_{(n1)}^1Y_{(n2)}X_{(n2)}^1Y_{(n3)}X_{(n3)}^1\mathrm{}Y_1X_1^1Y_0X_0^1\end{array}\right)X_0^n`$
where we have used the relation
$$Y_jX_k^1=X_{(k+1)}^1Y_j$$
due to $`a^{}`$ in $`Y_j`$. Moreover, by (148) and (149)
$$Y_jX_j^1=\sqrt{\frac{Nj}{N}}\frac{1}{R(Nj)+\theta }a^{}Z_j\text{for}j0.$$
Note that $`Z_0=Z`$. Therefore by using (147) we have
$$𝒜_n=\left(\begin{array}{c}\mathrm{𝟏}\\ \sqrt{{}_{n}{}^{}C_{1}^{}}Z_0\\ \mathrm{}\\ \sqrt{{}_{n}{}^{}C_{j}^{}}Z_{(j1)}Z_{(j2)}\mathrm{}Z_1Z_0\\ \mathrm{}\\ \sqrt{{}_{n}{}^{}C_{n1}^{}}Z_{(n2)}Z_{(n3)}\mathrm{}Z_1Z_0\\ Z_{(n1)}Z_{(n2)}Z_{(n3)}\mathrm{}Z_1Z_0\end{array}\right)\left(\mathrm{𝟏}+Z_0^{}Z_0\right)^{n/2}.$$
(156)
Now if we define
$$𝒵_n=\left(\begin{array}{c}\sqrt{{}_{n}{}^{}C_{1}^{}}Z_0\\ \mathrm{}\\ \sqrt{{}_{n}{}^{}C_{j}^{}}Z_{(j1)}Z_{(j2)}\mathrm{}Z_1Z_0\\ \mathrm{}\\ \sqrt{{}_{n}{}^{}C_{n1}^{}}Z_{(n2)}Z_{(n3)}\mathrm{}Z_1Z_0\\ Z_{(n1)}Z_{(n2)}Z_{(n3)}\mathrm{}Z_1Z_0\end{array}\right),$$
(157)
then
$$𝒜_n=\left(\begin{array}{c}\mathrm{𝟏}\\ 𝒵_n\end{array}\right)\left(\mathrm{𝟏}+Z_0^{}Z_0\right)^{n/2}.$$
and it is easy to show
$$\mathrm{𝟏}+𝒵_n^{}𝒵_n=\left(\mathrm{𝟏}+Z_0^{}Z_0\right)^n,$$
so we obtain
$`𝒫_n`$ $`=`$ $`𝒜_n𝒜_n^{}`$ (160)
$`=`$ $`\left(\begin{array}{c}\mathrm{𝟏}\\ 𝒵_n\end{array}\right)\left(\mathrm{𝟏}+𝒵_n^{}𝒵_n\right)^1(\mathrm{𝟏},𝒵_n^{})`$
$`=`$ $`\left(\begin{array}{cc}\left(\mathrm{𝟏}+𝒵_n^{}𝒵_n\right)^1& \left(\mathrm{𝟏}+𝒵_n^{}𝒵_n\right)^1𝒵_n^{}\\ 𝒵_n\left(\mathrm{𝟏}+𝒵_n^{}𝒵_n\right)^1& 𝒵_n\left(\mathrm{𝟏}+𝒵_n^{}𝒵_n\right)^1𝒵_n^{}\end{array}\right)`$ (163)
$`=`$ $`\left(\begin{array}{cc}\mathrm{𝟏}& 𝒵_n^{}\\ 𝒵_n& \mathrm{𝟏}\end{array}\right)\left(\begin{array}{cc}\mathrm{𝟏}& \\ & \mathrm{𝟎}\end{array}\right)\left(\begin{array}{cc}\mathrm{𝟏}& 𝒵_n^{}\\ 𝒵_n& \mathrm{𝟏}\end{array}\right)^1.`$ (170)
This is the Oike expression in , see also Appendix B.
A comment is in order. Two of important properties which the classical Veronese mapping has are
1. The Veronese mapping $`𝐂P^1𝐂P^n`$ has the mapping degree $`n`$
2. The Veronese surface (which is the image of Veronese mapping) is a minimal surface in $`𝐂P^n`$
Since we have constructed a non–commutative version of the Veronese mapping, a natural question arises : What are non–commutative versions corresponding to 1. and 2. above ?
These are very interesting problems from the view point of non–commutative “differential” geometry. It is worth challenging.
### 6.2 Non–Commutative Pseudo Veronese Mapping
We make a review of the noncompact one of Veronese mapping which we call a pseudo Veronese mapping. First let us define the manifold $`𝐂Q^n`$ which is not always well known.
$$𝐂Q^n=\left\{QM(n+1;𝐂)\right|Q^2=Q,J_nQJ_n=Q^{}\text{and}\text{tr}Q=1\}$$
(171)
where $`J_n`$ is a matrix defined by
$$J_n=\left(\begin{array}{ccccc}1& & & & \\ & 1& & & \\ & & & & \\ & & & & \\ & & & & (1)^n\end{array}\right)M(n+1;𝐂).$$
We note that this $`J_n`$ is not a convensional one. Usually it is taken as
$$\stackrel{~}{J}_n=\left(\begin{array}{ccccc}1& & & & \\ & 1& & & \\ & & & & \\ & & & & \\ & & & & 1\end{array}\right)M(n+1;𝐂).$$
For the space $`H_𝐂^n`$ defined by
$$H_𝐂^n=\left\{v𝐂^{n+1}\right|v^{}J_nv=1\}$$
(172)
we can define a map
$$H_𝐂^n𝐂Q^n:vQ=vv^{}J_n.$$
(173)
For $`Q_1𝐂Q^1`$ it can be written as
$`Q_1`$ $`=`$ $`\left(\begin{array}{cc}\alpha & \overline{\beta }\\ \beta & \overline{\alpha }\end{array}\right)\left(\begin{array}{ccccc}1& & & & \\ & 0& & & \end{array}\right)\left(\begin{array}{cc}\alpha & \overline{\beta }\\ \beta & \overline{\alpha }\end{array}\right)^1,|\alpha |^2|\beta |^2=1`$
$`=`$ $`v_1v_1^{}J_1,`$
where
$$v_1=\left(\begin{array}{c}\alpha \\ \beta \end{array}\right)H_𝐂^1.$$
For this $`v_1`$ we define $`v_n`$ as
$$v_n=\left(\begin{array}{c}\alpha ^n\\ \sqrt{{}_{n}{}^{}C_{1}^{}}\alpha ^{n1}(\beta )\\ \mathrm{}\\ \sqrt{{}_{n}{}^{}C_{j}^{}}\alpha ^{nj}(\beta )^j\\ \mathrm{}\\ (\beta )^n\end{array}\right).$$
(175)
Then it is easy to see
$$v_n^{}J_nv_n=\left(|\alpha |^2|\beta |^2\right)^n=1,$$
so $`v_nH_𝐂^n`$. Namely, we defined the map
$$H_𝐂^1H_𝐂^n:v_1v_n.$$
Therefore, we have the noncompact one of Veronese mapping
$$𝐂Q^1𝐂Q^n:Q_1=v_1v_1^{}J_1Q_n=v_nv_n^{}J_n.$$
(176)
For a tentative terminology let us call this a pseudo Veronese mapping.
Next, let us consider a local coordinate system. From (175)
$$v_n=\left(\begin{array}{c}1\\ \sqrt{{}_{n}{}^{}C_{1}^{}}w\\ \mathrm{}\\ \sqrt{{}_{n}{}^{}C_{j}^{}}w^j\\ \mathrm{}\\ w^n\end{array}\right)\alpha ^n\text{where}w=\alpha /\beta $$
then it is easy to check $`|w|^2<1`$. We define a domain like (open) hyperbolic pillar
$$D_J^n=\left\{v𝐂^n\right|v^{}J_{n1}v<1\}.$$
Then
$$D_J^1D_J^n:w\left(\begin{array}{c}\sqrt{{}_{n}{}^{}C_{1}^{}}w\\ \mathrm{}\\ \sqrt{{}_{n}{}^{}C_{j}^{}}w^j\\ \mathrm{}\\ w^n\end{array}\right)$$
is a local map that we are looking for. As a whole see the following picture.
Next we want to consider a non–commutative version of the map. If we set
$$\left(\begin{array}{c}\mathrm{\Gamma }_0\\ \mathrm{\Omega }_0\end{array}\right)=\left(\begin{array}{c}\frac{S(N+1)+\theta }{\sqrt{2S(N+1)(S(N+1)+\theta )}}\\ \frac{1}{\sqrt{2S(N)(S(N)+\theta )}}a^{}\end{array}\right)$$
(177)
from $`V`$ in (124), then
$$^{}𝐉=\mathrm{\Gamma }_0^2\mathrm{\Omega }_0^{}\mathrm{\Omega }_0=\mathrm{𝟏}\text{and}\mathrm{\Omega }_0\mathrm{\Gamma }_0^1=\frac{1}{S(N)+\theta }a^{}W.$$
That is, $`=(\mathrm{\Gamma }_0,\mathrm{\Omega }_0)^T`$ is a non–commutative hyperboloid and $`W`$ is a kind of “stereographic projection” of the hyperboloid. It is easy to see the following
$$\mathrm{𝟏}W^{}W=\frac{2S(N+1)}{S(N+1)+\theta }=\mathrm{\Gamma }_0^2\mathrm{\Gamma }_0=\left(\mathrm{𝟏}W^{}W\right)^{1/2}.$$
(178)
Here let us introduce new notations for the following. For $`j0`$ we set
$`\mathrm{\Gamma }_j`$ $`=`$ $`{\displaystyle \frac{S(N+1j)+\theta }{\sqrt{2S(N+1j)(S(N+1j)+\theta )}}},`$ (179)
$`\mathrm{\Omega }_j`$ $`=`$ $`\sqrt{{\displaystyle \frac{Nj}{N}}}{\displaystyle \frac{1}{\sqrt{2S(Nj)(S(Nj)+\theta )}}}a^{}.`$ (180)
We list some useful formulas.
$$\mathrm{\Gamma }_j^2\mathrm{\Omega }_j^{}\mathrm{\Omega }_j=\mathrm{𝟏}\text{and}\mathrm{\Omega }_j^{}\mathrm{\Omega }_j=\mathrm{\Omega }_{j+1}\mathrm{\Omega }_{j+1}^{}\text{for}j0.$$
(181)
Now we are in a position to define a non–commutative version of the pseudo Veronese mapping.
$$=\left(\begin{array}{c}\mathrm{\Gamma }_0\\ \mathrm{\Omega }_0\end{array}\right)_n=\left(\begin{array}{c}\mathrm{\Gamma }_0^n\\ \sqrt{{}_{n}{}^{}C_{1}^{}}\mathrm{\Omega }_0\mathrm{\Gamma }_0^{n1}\\ \mathrm{}\\ \sqrt{{}_{n}{}^{}C_{j}^{}}\mathrm{\Omega }_{(j1)}\mathrm{\Omega }_{(j2)}\mathrm{}\mathrm{\Omega }_1\mathrm{\Omega }_0\mathrm{\Gamma }_0^{nj}\\ \mathrm{}\\ \sqrt{{}_{n}{}^{}C_{n1}^{}}\mathrm{\Omega }_{(n2)}\mathrm{\Omega }_{(n3)}\mathrm{}\mathrm{\Omega }_1\mathrm{\Omega }_0\mathrm{\Gamma }_0\\ \mathrm{\Omega }_{(n1)}\mathrm{\Omega }_{(n2)}\mathrm{\Omega }_{(n3)}\mathrm{}\mathrm{\Omega }_1\mathrm{\Omega }_0\end{array}\right).$$
(182)
Then it is not difficult to see
$$_n^{}𝐉_n_n=\left(\mathrm{\Gamma }_0^2\mathrm{\Omega }_0^{}\mathrm{\Omega }_0\right)^n=\left(^{}𝐉\right)^n=\mathrm{𝟏},$$
where $`𝐉_n`$ ($`𝐉_1=𝐉`$) is defined by
$$𝐉_n=\left(\begin{array}{ccccc}\mathrm{𝟏}& & & & \\ & \mathrm{𝟏}& & & \\ & & & & \\ & & & & \\ & & & & (1)^n\mathrm{𝟏}\end{array}\right)=J_n\mathrm{𝟏}.$$
From this we can define the projectors which correspond to pseudo projective spaces like
$$𝒬_n=_n_n^{}𝐉_n,𝒬_1=^{}𝐉,$$
(183)
so the map
$$𝒬_1𝒬_n$$
(184)
is a non-commutative version of the pseudo Veronese mapping.
## 7 Non–Commutative Representation Theory
In this section we construct a map (in the non–commutative models) corresponding to spin $`j`$–representation for the compact group $`SU(2)`$ and noncompact group $`SU(1,1)`$$`\mathrm{}`$ a kind of non–commutative version of classical spin representations $`\mathrm{}`$.
### 7.1 Non–Commutative Version of $`SU(2)`$ Case
The construction of spin $`j`$–representation ($`j𝐙_0+1/2`$) is well–known. Let us make a brief review within our necessity. For the vector space
$$_J=\text{Vect}_𝐂\left\{\sqrt{{}_{J1}{}^{}C_{k}^{}}z^k\right|k\{0,1,\mathrm{},J1\}\}$$
where $`J=2j+1(𝐍)`$, the inner product in this space is given by
$$<f|g>=\frac{2J}{2\pi }_𝐂\frac{d^2z}{(1+|z|^2)^{J+1}}f(z)\overline{g(z)}=\underset{k=0}{\overset{J1}{}}a_k\overline{b}_k$$
for $`f(z)=_{k=0}^{J1}\sqrt{{}_{J1}{}^{}C_{k}^{}}a_kz^k`$ and $`g(z)=_{k=0}^{J1}\sqrt{{}_{J1}{}^{}C_{k}^{}}b_kz^k`$ in $`_J`$. Here $`d^2z`$ means $`dxdy`$ for $`z=x+iy`$.
For example, for $`j=1/2`$, $`j=1`$ and $`j=3/2`$
$$_2=\text{Vect}_𝐂\{1,z\},_3=\text{Vect}_𝐂\{1,\sqrt{2}z,z^2\},_4=\text{Vect}_𝐂\{1,\sqrt{3}z,\sqrt{3}z^2,z^3\}.$$
Therefore, we identify $`_J`$ with $`𝐂^J`$ by
$$f(z)=\underset{k=0}{\overset{J1}{}}\sqrt{{}_{J1}{}^{}C_{k}^{}}a_kz^k(a_0,a_1,\mathrm{},a_{J1})^\text{T}.$$
For
$$A=\left(\begin{array}{cc}\alpha & \overline{\beta }\\ \beta & \overline{\alpha }\end{array}\right)SU(2)\left(|\alpha |^2+|\beta |^2=1\right)$$
the spin $`j`$ representation
$$\varphi _j:SU(2)SU(J)$$
is defined as
$$\left(\varphi _j(A)f\right)(z)=(\alpha +\beta z)^{J1}f\left(\frac{\overline{\beta }+\overline{\alpha }z}{\alpha +\beta z}\right)$$
(185)
where $`f_J`$. It is easy to obtain $`\varphi _j(A)`$ for $`j=1/2`$, $`j=1`$ and $`j=3/2`$.
Namely, the spin 1/2 representation is
$$\varphi _{1/2}(A)=\left(\begin{array}{cc}\alpha & \overline{\beta }\\ \beta & \overline{\alpha }\end{array}\right)=A,$$
(186)
the spin 1 representation is
$$\varphi _1(A)=\left(\begin{array}{ccc}\alpha ^2& \sqrt{2}\alpha \overline{\beta }& \overline{\beta }^2\\ \sqrt{2}\alpha \beta & |\alpha |^2|\beta |^2& \sqrt{2}\overline{\alpha }\overline{\beta }\\ \beta ^2& \sqrt{2}\overline{\alpha }\beta & \overline{\alpha }^2\end{array}\right),$$
(187)
and the spin 3/2 representation is
$$\varphi _{3/2}(A)=\left(\begin{array}{cccc}\alpha ^3& \sqrt{3}\alpha ^2\overline{\beta }& \sqrt{3}\alpha \overline{\beta }^2& \overline{\beta }^3\\ \sqrt{3}\alpha ^2\beta & (|\alpha |^22|\beta |^2)\alpha & (2|\alpha |^2|\beta |^2)\overline{\beta }& \sqrt{3}\overline{\alpha }\overline{\beta }^2\\ \sqrt{3}\alpha \beta ^2& (2|\alpha |^2|\beta |^2)\beta & (|\alpha |^22|\beta |^2)\overline{\alpha }& \sqrt{3}\overline{\alpha }^2\overline{\beta }\\ \beta ^3& \sqrt{3}\overline{\alpha }\beta ^2& \sqrt{3}\overline{\alpha }^2\beta & \overline{\alpha }^3\end{array}\right).$$
(188)
Next we want to consider a non–commutative version of the spin representation. However, since such a theory has not been known as far as we know we must look for mappings corresponding to $`\varphi _1(A)`$ and $`\varphi _{3/2}(A)`$ by (many) trial and error, see Appendix C.
If we set
$$UU_I=\left(\begin{array}{cc}X_0& Y_0^{}\\ Y_0& X_1\end{array}\right):\text{unitary}$$
from (77), then the corresponding map for $`\varphi _1(A)`$ is
$$\mathrm{\Phi }_1(U)=\left(\begin{array}{ccc}X_0^2& \sqrt{2}X_0Y_0^{}& Y_0^{}Y_1^{}\\ \sqrt{2}Y_0X_0& X_1^2Y_1^{}Y_1& \sqrt{2}X_1Y_1^{}\\ Y_1Y_0& \sqrt{2}Y_1X_1& X_2^2\end{array}\right)$$
(189)
and the corresponding map for $`\varphi _{3/2}(A)`$ is
$`\mathrm{\Phi }_{3/2}(U)`$ $`=`$ $`\left(\begin{array}{cccc}X_0^3& \sqrt{3}X_0^2Y_0^{}& \sqrt{3}X_0Y_0^{}Y_1^{}& Y_0^{}Y_1^{}Y_2^{}\\ \sqrt{3}Y_0X_0^2& X_1\left(X_1^22Y_1^{}Y_1\right)& \left(2X_1^2Y_1^{}Y_1\right)Y_1^{}& \sqrt{3}X_1Y_1^{}Y_2^{}\\ \sqrt{3}Y_1Y_0X_0& Y_1\left(2X_1^2Y_1^{}Y_1\right)& X_2\left(X_2^22Y_2^{}Y_2\right)& \sqrt{3}X_2^2Y_2^{}\\ Y_2Y_1Y_0& \sqrt{3}Y_2Y_1X_1& \sqrt{3}Y_2X_2^2& X_3^3\end{array}\right).`$ (194)
To check the unitarity of $`\mathrm{\Phi }_1(U)`$ and $`\mathrm{\Phi }_{3/2}(U)`$ is long but straightforward.
For $`j2`$ we could not find a general method like (185) which determines $`\mathrm{\Phi }_j(U)`$. However, we know only that the first column of $`\mathrm{\Phi }_j(U)`$ is just $`𝒜_{2j}`$ in (151).,
$$\mathrm{\Phi }_j(U)=(𝒜_{2j},,\mathrm{},):\text{unitary}$$
and
$$\mathrm{\Phi }_j(U)\left(\begin{array}{ccccc}\mathrm{𝟏}& & & & \\ & \mathrm{𝟎}& & & \\ & & & & \\ & & & & \\ & & & & \mathrm{𝟎}\end{array}\right)\mathrm{\Phi }_j(U)^{}=𝒜_{2j}𝒜_{2j}^{}=𝒫_{2j}.$$
We leave finding a general method to the readers as a challenging problem.
### 7.2 Non–Commutative Version of $`SU(1,1)`$ Case
Let us review some aspects of the theory of unitary representation of $`SU(1,1)`$ within our necessity.
Let $`H^2H^2(D)`$ be the second Hardy class where $`D`$ is the open unit disk in $`𝐂`$. We consider the spin $`j`$ representation of the non–compact group $`SU(1,1)`$. The inner product is defined as
$$<f|g>=\frac{2(2j1)}{2\pi }_Dd^2z(1|z|^2)^{2j2}f(z)\overline{g(z)}=\underset{n=0}{\overset{\mathrm{}}{}}\frac{n!}{(2j)_n}a_n\overline{b}_n$$
for $`f(z)=_{n=0}^{\mathrm{}}a_nz^n`$ and $`g(z)=_{n=0}^{\mathrm{}}b_nz^n`$ in $`H^2`$. For $`j=\frac{1}{2}`$ we must take some renormalization into consideration (we omit it here). Then $`\{H^2,<|>\}`$ becomes a complex Hibert space.
Therefore, it is better for us to consider the vector space
$$H_{2j}^2=\text{Vect}_𝐂\{1,\sqrt{2j}z,\mathrm{},\sqrt{\frac{(2j)_k}{k!}}z^k,\mathrm{}\}$$
and the correspondence between $`H_{2j}^2`$ and $`\mathrm{}^2(𝐂)`$ is given by
$$f(z)=\underset{n=0}{\overset{\mathrm{}}{}}\sqrt{\frac{(2j)_n}{n!}}a_nz^n(a_0,a_1,\mathrm{},a_n,\mathrm{})^\text{T},$$
so we identify $`H_{2j}^2`$ with $`\mathrm{}^2(𝐂)`$ by this correspondence.
For
$$B=\left(\begin{array}{cc}\alpha & \overline{\beta }\\ \beta & \overline{\alpha }\end{array}\right)SU(1,1)\left(|\alpha |^2|\beta |^2=1\right)$$
the spin $`j`$ unitary representation
$$\psi _j:SU(1,1)U(\mathrm{}^2(𝐂))$$
is defined as
$$\left(\psi _j(B)f\right)(z)=(\alpha +\beta z)^{2j}f\left(\frac{\overline{\beta }+\overline{\alpha }z}{\alpha +\beta z}\right)$$
(196)
where $`fH_{2j}^2`$.
For example, when $`f=1`$ (a constant) it is easy to see
$`\left(\psi _j(B)1\right)(z)`$ $`=`$ $`(\alpha +\beta z)^{2j}`$
$`=`$ $`{\displaystyle \frac{1}{\alpha ^{2j}}}2j{\displaystyle \frac{\beta }{\alpha ^{2j+1}}}z+\mathrm{}+(1)^n{\displaystyle \frac{(2j)_n}{n!}}{\displaystyle \frac{\beta ^n}{\alpha ^{2j+n}}}z^n+\mathrm{}`$
$`=`$ $`{\displaystyle \frac{1}{\alpha ^{2j}}}\sqrt{2j}{\displaystyle \frac{\beta }{\alpha ^{2j+1}}}\sqrt{2j}z+\mathrm{}+(1)^n\sqrt{{\displaystyle \frac{(2j)_n}{n!}}}{\displaystyle \frac{\beta ^n}{\alpha ^{2j+n}}}\sqrt{{\displaystyle \frac{(2j)_n}{n!}}}z^n+\mathrm{}`$
where $`(a)_n`$ is the Pochammer notation defined by
$$(a)_n=a(a+1)\mathrm{}(a+n1).$$
Therefore
$$\psi _j(B)1=\left(\begin{array}{c}\frac{1}{\alpha ^{2j}}\\ \sqrt{2j}\frac{\beta }{\alpha ^{2j+1}}\\ \mathrm{}\\ (1)^n\sqrt{\frac{(2j)_n}{n!}}\frac{\beta ^n}{\alpha ^{2j+n}}\\ \mathrm{}\end{array}\right)=\left(\begin{array}{c}\alpha ^{2j}\\ \sqrt{2j}\beta \alpha ^{(2j+1)}\\ \mathrm{}\\ (1)^n\sqrt{\frac{(2j)_n}{n!}}\beta ^n\alpha ^{(2j+n)}\\ \mathrm{}\end{array}\right).$$
(197)
More generally, for $`f_k(z)=\sqrt{\frac{(2j)_k}{k!}}z^k`$
$`\left(\psi _j(B)f_k\right)(z)`$ $`=`$ $`\sqrt{{\displaystyle \frac{(2j)_k}{k!}}}(\overline{\beta }+\overline{\alpha }z)^k(\alpha +\beta z)^{(2j+k)}`$
$`=`$ $`\sqrt{{\displaystyle \frac{(2j)_k}{k!}}}{\displaystyle \underset{l=0}{\overset{k}{}}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{}_{k}{}^{}C_{l}^{}(1)^n{\displaystyle \frac{(2j+k)_n}{n!}}{\displaystyle \frac{\beta ^n\overline{\beta }^{kl}\overline{\alpha }^l}{\alpha ^{2j+n+k}}}z^{l+n},`$
so
$$\psi _j(B)f_k=\left(\begin{array}{c}\sqrt{\frac{(2j)_k}{k!}}\frac{\overline{\beta }^k}{\alpha ^{2j+k}}\\ \sqrt{\frac{(2j)_k}{k!}}\sqrt{\frac{1}{2j}}\left\{k|\alpha |^2(2j+k)|\beta |^2\right\}\frac{\overline{\beta }^{k1}}{\alpha ^{2j+k+1}}\\ \sqrt{\frac{(2j)_k}{k!}}\sqrt{\frac{2!}{(2j)_2}}\left\{\frac{k(k1)}{2}|\alpha |^4k(2j+k)|\alpha |^2|\beta |^2+\frac{(2j+k)(2j+k+1)}{2}|\beta |^4\right\}\frac{\overline{\beta }^{k2}}{\alpha ^{2j+k+2}}\\ \mathrm{}\\ \mathrm{}\end{array}\right).$$
Therefore the matrix defined by
$$\psi _j(B)=(\psi _j(B)f_0,\psi _j(B)f_1,\mathrm{},\psi _j(B)f_k,\mathrm{})$$
(198)
is the unitary representation that we are looking for.
We set
$$V=\left(\begin{array}{cc}\mathrm{\Gamma }_0& \mathrm{\Omega }_0^{}\\ \mathrm{\Omega }_0& \mathrm{\Gamma }_1\end{array}\right):\text{pseudo unitary}(V^{}𝐉V=𝐉)$$
from (124). In this case it is almost impossible to obtain the explicit unitary operator $`\mathrm{\Psi }_j(V)`$ corresponding to $`\psi _j(B)`$.
However, we can at least determine the first column of $`\mathrm{\Psi }_j(V)`$ by making use of (197) :
$$\widehat{}\left(\begin{array}{c}\mathrm{\Gamma }_0^{2j}\\ \sqrt{2j}\mathrm{\Omega }_0\mathrm{\Gamma }_0^{(2j+1)}\\ \mathrm{}\\ \sqrt{\frac{(2j)_n}{n!}}\mathrm{\Omega }_{(n1)}\mathrm{\Omega }_{(n2)}\mathrm{\Omega }_{(n3)}\mathrm{}\mathrm{\Omega }_1\mathrm{\Omega }_0\mathrm{\Gamma }_0^{(2j+n)}\\ \mathrm{}\end{array}\right)$$
(199)
with $`\mathrm{\Gamma }_0`$ and $`\mathrm{\Omega }_j`$ in (179) and (180). Then it is not difficult to see
$$\widehat{}^{}\widehat{}=\left(\mathrm{\Gamma }_0^2\mathrm{\Omega }_0^{}\mathrm{\Omega }_0\right)^{2j}=\mathrm{𝟏}.$$
Therefore, the map making use of projectors
$$𝒬_1=^{}𝐉\widehat{𝒫}=\widehat{}\widehat{}^{}$$
(200)
is a non–commutative version of the unitary expression of pseudo Veronese mapping. Compare the discussion here with the one after the equation (182).
We note that if the unitary operator
$$\mathrm{\Psi }_j(V)=(\widehat{}_0,\widehat{}_1,\mathrm{},\widehat{}_n,\mathrm{}),\widehat{}_0=\widehat{}$$
could be defined (we cannot determine $`\widehat{}_n`$ for $`n1`$), then we have
$$\mathrm{\Psi }_j(V)\left(\begin{array}{cccc}\mathrm{𝟏}& & & \\ & \mathrm{𝟎}& & \\ & & \mathrm{𝟎}& \\ & & & \mathrm{}\end{array}\right)\mathrm{\Psi }_j(V)^{}=\widehat{}\widehat{}^{}=\widehat{𝒫}.$$
A comment is in order. In the construction of $`\widehat{}_n`$ we need an infinite number of operators, which means a kind of second non–commutativization.
## 8 Discussion
In this paper we derived a non–commutative version of the Berry model (based on $`SU(2)`$) arising from the Jaynes–Cummings model in quantum optics and the pseudo Berry model (based on $`SU(1,1)`$) by changing the generators, and constructed a non–commutative version of the Hopf and pseudo Hopf bundles in the classical case.
The bundle has a kind of Dirac strings in the case of non–commutative Berry model. However, they appear in only states containing the ground one ($`\times \{|0\}\{|0\}\times \times `$) and don’t appear in excited states, which is very interesting.
In general, a non-commutative version of classical field theory is of course not unique. If our model is a “correct” one, then this paper give an example that classical singularities like Dirac strings are not universal in some non–commutative model. As to general case with higher spins which are not easy, see .
Moreover, for the two models a non–commutative version of the Veronese mapping or pseudo Veronese mapping was constructed, and unitary mappings corresponding to (classical) spin representations were constructed though they are not necessarily enough.
The results or methods in the paper will become a starting point to construct a fruitful non–commutative geometry or representation theory.
Last, we would like to make a comment. To develop a “quantum” mathematics we need a rigorous method to treat an analysis or a geometry on infinite dimensional spaces like Fock space. In quantum field theories physicists have given some (interesting) methods, while they are more or less formal from the mathematical point of view. It is a rigorous method which we need. As a trial is recommended.
Acknowledgment.
The author wishes to thank Akira Asada, Yoshinori Machida, Shin’ichi Nojiri, Ryu Sasaki and Tatsuo Suzuki for their helpful comments and suggestions.
The author also thanks to Gennadi Sardanashvily and Giovanni Giachetta for warm hospitality at Firenze (14-18/April/2005). The arrangement of this paper was determined during the stay.
Appendix
A Classical Theory of Projective Spaces
Complex projective spaces are typical examples of symmetric spaces and are very tractable, so they are used to construct several examples in both physics and mathematics. We make a review of complex projective spaces within our necessity, see for example , , .
For $`n𝐍`$ the complex projective space $`𝐂P^n`$ is defined as follows : For $`𝜻`$, $`𝝁`$ $`𝐂^{n+1}\{\mathrm{𝟎}\}`$ $`𝜻`$ is equivalent to $`𝝁`$ ($`𝜻`$ $``$ $`𝝁`$) if and only if $`𝜻`$ = $`\lambda `$$`𝝁`$ for $`\lambda 𝐂\{0\}`$. We show the equivalent relation class as \[$`𝜻`$\] and set $`𝐂P^n𝐂^{n+1}\{\mathrm{𝟎}\}/`$. For $`𝜻`$ = $`(\zeta _0,\zeta _1,\mathrm{},\zeta _n)`$ we write usually as \[$`𝜻`$\] = $`[\zeta _0:\zeta _1:\mathrm{}:\zeta _n]`$. Then it is well–known that $`𝐂P^n`$ has $`n+1`$ local charts, namely
$$𝐂P^n=\underset{j=0}{\overset{n}{}}U_j,U_j=\{[\zeta _0:\mathrm{}:\zeta _j:\mathrm{}:\zeta _n]|\zeta _j0\}.$$
(201)
Since
$$(\zeta _0,\mathrm{},\zeta _j,\mathrm{},\zeta _n)=\zeta _j(\frac{\zeta _0}{\zeta _j},\mathrm{},\frac{\zeta _{j1}}{\zeta _j},1,\frac{\zeta _{j+1}}{\zeta _j},\mathrm{},\frac{\zeta _n}{\zeta _j}),$$
we have the local coordinate on $`U_j`$
$$(\frac{\zeta _0}{\zeta _j},\mathrm{},\frac{\zeta _{j1}}{\zeta _j},\frac{\zeta _{j+1}}{\zeta _j},\mathrm{},\frac{\zeta _n}{\zeta _j}).$$
(202)
However the above definition of $`𝐂P^n`$ is not tractable, so we use the well–known expression by projections (see )
$$𝐂P^nG_{1,n+1}(𝐂)=\{PM(n+1;𝐂)|P^2=P,P=P^{}\text{and}\text{tr}P=1\}$$
(203)
and the correspondence
$$[\zeta _0:\zeta _1:\mathrm{}:\zeta _n]\frac{1}{|\zeta _0|^2+|\zeta _1|^2+\mathrm{}+|\zeta _n|^2}\left(\begin{array}{ccccc}|\zeta _0|^2& \zeta _0\overline{\zeta }_1& & & \zeta _0\overline{\zeta }_n\\ \zeta _1\overline{\zeta }_0& |\zeta _1|^2& & & \zeta _1\overline{\zeta }_n\\ & & & & \\ & & & & \\ \zeta _n\overline{\zeta }_0& \zeta _n\overline{\zeta }_1& & & |\zeta _n|^2\end{array}\right)P.$$
(204)
If we set
$$|𝜻=\frac{1}{\sqrt{_{j=0}^n|\zeta _j|^2}}\left(\begin{array}{c}\zeta _0\\ \zeta _1\\ \\ \\ \zeta _n\end{array}\right),$$
(205)
then we can write the right hand side of (204) as
$$P=|𝜻𝜻|\text{and}𝜻|𝜻=1.$$
(206)
For example on $`U_0`$
$$(z_1,z_2,\mathrm{},z_n)=(\frac{\zeta _1}{\zeta _0},\frac{\zeta _2}{\zeta _0},\mathrm{},\frac{\zeta _n}{\zeta _0}),$$
we have
$`P(z_1,\mathrm{},z_n)`$ $`=`$ $`{\displaystyle \frac{1}{1+_{j=1}^n|z_j|^2}}\left(\begin{array}{ccccc}1& \overline{z}_1& & & \overline{z}_n\\ z_1& |z_1|^2& & & z_1\overline{z}_n\\ & & & & \\ & & & & \\ z_n& z_n\overline{z}_1& & & |z_n|^2\end{array}\right)`$ (212)
$`=`$ $`|(z_1,z_2,\mathrm{},z_n)(z_1,z_2,\mathrm{},z_n)|,`$ (213)
where
$$|(z_1,z_2,\mathrm{},z_n)=\frac{1}{\sqrt{1+_{j=1}^n|z_j|^2}}\left(\begin{array}{c}1\\ z_1\\ \\ \\ z_n\end{array}\right).$$
To be clearer, let us give a detailed description for the case of $`n`$ = $`1`$ and $`2`$.
(a) $`n=1`$ :
$`P(z)`$ $`=`$ $`{\displaystyle \frac{1}{1+|z|^2}}\left(\begin{array}{cc}1& \overline{z}\\ z& |z|^2\end{array}\right)=|zz|,`$ (219)
$`\text{where}|z={\displaystyle \frac{1}{\sqrt{1+|z|^2}}}\left(\begin{array}{c}1\\ z\end{array}\right),z={\displaystyle \frac{\zeta _1}{\zeta _0}},\text{on}U_0,`$
$`P(w)`$ $`=`$ $`{\displaystyle \frac{1}{|w|^2+1}}\left(\begin{array}{cc}|w|^2& w\\ \overline{w}& 1\end{array}\right)=|ww|,`$ (225)
$`\text{where}|w={\displaystyle \frac{1}{\sqrt{|w|^2+1}}}\left(\begin{array}{c}w\\ 1\end{array}\right),w={\displaystyle \frac{\zeta _0}{\zeta _1}},\text{on}U_1.`$
(b) $`n=2`$ :
$`P(z_1,z_2)`$ $`=`$ $`{\displaystyle \frac{1}{1+|z_1|^2+|z_2|^2}}\left(\begin{array}{ccc}1& \overline{z}_1& \overline{z}_2\\ z_1& |z_1|^2& z_1\overline{z}_2\\ z_2& z_2\overline{z}_1& |z_2|^2\end{array}\right)=|(z_1,z_2)(z_1,z_2)|,`$ (229)
where $`|(z_1,z_2)={\displaystyle \frac{1}{\sqrt{1+|z_1|^2+|z_2|^2}}}\left(\begin{array}{c}1\\ z_1\\ z_2\end{array}\right),(z_1,z_2)=({\displaystyle \frac{\zeta _1}{\zeta _0}},{\displaystyle \frac{\zeta _2}{\zeta _0}})\text{on}U_0,`$ (233)
$`P(w_1,w_2)`$ $`=`$ $`{\displaystyle \frac{1}{|w_1|^2+1+|w_2|^2}}\left(\begin{array}{ccc}|w_1|^2& w_1& w_1\overline{w}_2\\ \overline{w}_1& 1& \overline{w}_2\\ w_2\overline{w}_1& w_2& |w_2|^2\end{array}\right)=|(w_1,w_2)(w_1,w_2)|,`$ (237)
where $`|(w_1,w_2)={\displaystyle \frac{1}{\sqrt{|w_1|^2+1+|w_2|^2}}}\left(\begin{array}{c}w_1\\ 1\\ w_2\end{array}\right),(w_1,w_2)=({\displaystyle \frac{\zeta _0}{\zeta _1}},{\displaystyle \frac{\zeta _2}{\zeta _1}})\text{on}U_1,`$ (241)
$`P(v_1,v_2)`$ $`=`$ $`{\displaystyle \frac{1}{|v_1|^2+|v_2|^2+1}}\left(\begin{array}{ccc}|v_1|^2& v_1\overline{v}_2& v_1\\ v_2\overline{v}_1& |v_2|^2& v_2\\ \overline{v}_1& \overline{v}_2& 1\end{array}\right)=|(v_1,v_2)(v_1,v_2)|,`$ (245)
where $`|(v_1,v_2)={\displaystyle \frac{1}{\sqrt{|v_1|^2+|v_2|^2+1}}}\left(\begin{array}{c}v_1\\ v_2\\ 1\end{array}\right),(v_1,v_2)=({\displaystyle \frac{\zeta _0}{\zeta _2}},{\displaystyle \frac{\zeta _1}{\zeta _2}})\text{on}U_2.`$ (249)
B Local Coordinate of the Projector
We give a proof to the last formula in (160).
By making use of the expression by Oike in (we don’t repeat it here)
$$𝒫(𝒵)=\left(\begin{array}{cc}\mathrm{𝟏}& 𝒵^{}\\ 𝒵& \mathrm{𝟏}\end{array}\right)\left(\begin{array}{cc}\mathrm{𝟏}& \\ & \mathrm{𝟎}\end{array}\right)\left(\begin{array}{cc}\mathrm{𝟏}& 𝒵^{}\\ 𝒵& \mathrm{𝟏}\end{array}\right)^1$$
(250)
where $`𝒵`$ is some operator on the Fock space $``$. Let us rewrite this into more useful form. From the simple relation
$$\left(\begin{array}{cc}\mathrm{𝟏}& 𝒵^{}\\ 𝒵& \mathrm{𝟏}\end{array}\right)\left(\begin{array}{cc}\mathrm{𝟏}& 𝒵^{}\\ 𝒵& \mathrm{𝟏}\end{array}\right)=\left(\begin{array}{cc}\mathrm{𝟏}+𝒵^{}𝒵& \\ & \mathrm{𝟏}+𝒵𝒵^{}\end{array}\right)$$
we have
$$\left(\begin{array}{cc}\mathrm{𝟏}& 𝒵^{}\\ 𝒵& \mathrm{𝟏}\end{array}\right)^1=\left(\begin{array}{cc}(\mathrm{𝟏}+𝒵^{}𝒵)^1& \\ & (\mathrm{𝟏}+𝒵𝒵^{})^1\end{array}\right)\left(\begin{array}{cc}\mathrm{𝟏}& 𝒵^{}\\ 𝒵& \mathrm{𝟏}\end{array}\right).$$
Inserting this into (250) and some calculation leads to
$$𝒫(𝒵)=\left(\begin{array}{cc}(\mathrm{𝟏}+𝒵^{}𝒵)^1& (\mathrm{𝟏}+𝒵^{}𝒵)^1𝒵^{}\\ 𝒵(\mathrm{𝟏}+𝒵^{}𝒵)^1& 𝒵(\mathrm{𝟏}+𝒵^{}𝒵)^1𝒵^{}\end{array}\right).$$
(251)
Comparing (251) with (102) we obtain the “local coordinate”
$$𝒵=\frac{1}{R(N)+\theta }a^{}=a^{}\frac{1}{R(N+1)+\theta }$$
(252)
where $`R(N)=\sqrt{N+\theta ^2}`$. $`𝒵`$ obtained by “stereographic projection” is a kind of complex coordinate.
Now if we take a classical limit $`axiy`$, $`a^{}x+iy`$ and $`\theta =z`$ then
$$Z_c=\frac{x+iy}{r+z}$$
(253)
where $`r=\sqrt{x^2+y^2+z^2}`$. This is nothing but a well–known one for (16).
C Some Calculations of First Chern Class
We calculate the first Chern class of some vector bundles on $`𝐂P^1`$ and show that the mapping degree of Veronese mapping is just $`n`$.
We write our definition of $`𝐂P^n`$ once more :
$$𝐂P^n=\{PM(n+1;𝐂)|P^2=P,P=P^{}\text{and}\text{tr}P=1\}.$$
On this space we define a canonical vector bundle like
$`E_n`$ $`=`$ $`\left\{(P,v)𝐂P^n\times 𝐂^{n+1}\right|Pv=v\},`$
$`\pi `$ $`:`$ $`E_n𝐂P^n,\pi (P,v)=P.`$
Then the system $`\xi _n=\{𝐂,E_n,\pi ,𝐂P^n\}`$ is called the canonical line bundle (because $`P`$ is rank one), see , . This is one of most important vector bundles.
Let us calculate the first Chern class of $`\xi _1`$. For the local coordinate $`z`$ in section 6.1, $`P`$ can be written as
$$P(z)=\frac{1}{1+|z|^2}\left(\begin{array}{cc}1& \overline{z}\\ z& |z|^2\end{array}\right),v(z)=\alpha \left(\begin{array}{c}1\\ z\end{array}\right)(\alpha 𝐂).$$
(254)
Then the canonical connection $`𝒜`$ and its curvature $``$ can be written as
$$𝒜=\frac{\overline{z}}{1+|z|^2}dz,=d𝒜=\frac{1}{(1+|z|^2)^2}d\overline{z}dz.$$
(255)
Let $`\chi `$ be the Veronese mapping in section 6.1 ($`\chi :𝐂P^1𝐂P^n`$). then we can consider the pull–back bundle $`\chi ^{}\xi _n=\{𝐂,\chi ^{}(E_n),\pi ,𝐂P^1\}`$ where
$`\chi ^{}(E_n)`$ $`=`$ $`\left\{(P,v)𝐂P^1\times 𝐂^{n+1}\right|\chi (P)v=v\}`$
$`\pi `$ $`:`$ $`\chi ^{}(E_n)𝐂P^1,\pi (P,v)=P.`$
See the following picture.
Let us give a local description. For $`z`$ in (254)
$$\chi (P(z))=\frac{1}{(1+|z|^2)^n}\left(\begin{array}{cc}1& \psi (z)^{}\\ \psi (z)& \psi (z)\psi (z)^{}\end{array}\right),v(z)=\alpha \left(\begin{array}{c}1\\ \psi (z)\end{array}\right)(\alpha 𝐂)$$
where $`\psi (z)`$ is the map defined in section 6.1
$$\psi (z)=\left(\begin{array}{c}\sqrt{{}_{n}{}^{}C_{1}^{}}z\\ \mathrm{}\\ \sqrt{{}_{n}{}^{}C_{j}^{}}z^j\\ \mathrm{}\\ \sqrt{{}_{n}{}^{}C_{n1}^{}}z^{n1}\\ z^n\end{array}\right)1+\psi (z)^{}\psi (z)=(1+|z|^2)^n.$$
Now the connection and curvature of the pull–backed bundle are given by
$$𝒜_n=(1+\psi (z)^{}\psi (z))^1\psi (z)^{}d\psi (z),_n=d𝒜_n.$$
(256)
Let us calculate : it is easy to see
$`𝒜_n`$ $`=`$ $`{\displaystyle \frac{{}_{n}{}^{}C_{1}^{}+\mathrm{}+j{}_{n}{}^{}C_{j}^{}|z|^{2(j1)}+\mathrm{}+n{}_{n}{}^{}C_{n}^{}|z|^{2(n1)}}{(1+|z|^2)^n}}\overline{z}dz`$
$`=`$ $`{\displaystyle \frac{\frac{d}{d(|z|^2)}({}_{n}{}^{}C_{1}^{}|z|^2++\mathrm{}+{}_{n}{}^{}C_{j}^{}|z|^{2j}+\mathrm{}+{}_{n}{}^{}C_{n}^{}|z|^{2n})}{(1+|z|^2)^n}}\overline{z}dz`$
$`=`$ $`{\displaystyle \frac{\frac{d}{d(|z|^2)}\left((1+|z|^2)^n1\right)}{(1+|z|^2)^n}}\overline{z}dz`$
$`=`$ $`{\displaystyle \frac{n(1+|z|^2)^{n1}}{(1+|z|^2)^n}}\overline{z}dz`$
$`=`$ $`n{\displaystyle \frac{\overline{z}}{1+|z|^2}}dz`$
$`=`$ $`n𝒜,`$
therefore
$$_n=n=n\frac{1}{(1+|z|^2)^2}d\overline{z}dz.$$
As a result we have
$$\text{Ch}_1(\chi ^{}\xi _n)=\frac{1}{2\pi i}_𝐂n\frac{1}{(1+|z|^2)^2}𝑑\overline{z}dz=n.$$
(257)
As to calculations of geometric objects like Chern classes or holonomies on quantum computation see for example or .
D Difficulty of Tensor Decomposition
We point out a difficulty in obtaining the formula (189) or (194) by decomposing tensor products of $`V`$.
To obtain the formula (187) there is another method which uses a decomposition of the tensor product $`AA`$. Let us introduce. For
$$A=\left(\begin{array}{cc}\alpha & \overline{\beta }\\ \beta & \overline{\alpha }\end{array}\right)SU(2)$$
we have
$$AA=\left(\begin{array}{cccc}\alpha ^2& \alpha \overline{\beta }& \alpha \overline{\beta }& \overline{\beta }^2\\ \alpha \beta & |\alpha |^2& |\beta |^2& \overline{\alpha }\overline{\beta }\\ \alpha \beta & |\beta |^2& |\alpha |^2& \overline{\alpha }\overline{\beta }\\ \beta ^2& \overline{\alpha }\beta & \overline{\alpha }\beta & \overline{\alpha }^2\end{array}\right).$$
For the matrix $`T`$ coming from the Clebsch–Gordan decomposition
$$T=\left(\begin{array}{cccc}0& 1& 0& 0\\ \frac{1}{\sqrt{2}}& 0& \frac{1}{\sqrt{2}}& 0\\ \frac{1}{\sqrt{2}}& 0& \frac{1}{\sqrt{2}}& 0\\ 0& 0& 0& 1\end{array}\right)$$
it is easy to see
$$T^{}(AA)T=\left(\begin{array}{cccc}|\alpha |^2+|\beta |^2& & & \\ & \alpha ^2& \sqrt{2}\alpha \overline{\beta }& \overline{\beta }^2\\ & \sqrt{2}\alpha \beta & |\alpha |^2|\beta |^2& \sqrt{2}\overline{\alpha }\overline{\beta }\\ & \beta ^2& \sqrt{2}\overline{\alpha }\beta & \overline{\alpha }^2\end{array}\right)=\left(\begin{array}{cc}1& \\ & \varphi _1(A)\end{array}\right)$$
(258)
where we have used $`|\alpha |^2+|\beta |^2=1`$. This means a well–known decomposition
$$\frac{1}{2}\frac{1}{2}=01.$$
Let us take an analogy. For
$$V=\left(\begin{array}{cc}X_0& Y_0^{}\\ Y_0& X_1\end{array}\right)$$
we have
$$VV=\left(\begin{array}{cccc}X_0^2& X_0Y_0^{}& Y_0^{}X_0& Y_0^{}Y_0^{}\\ X_0Y_0& X_0X_1& Y_0^{}Y_0& Y_0^{}X_1\\ Y_0X_0& Y_0Y_0^{}& X_1X_0& X_1Y_0^{}\\ Y_0Y_0& Y_0X_1& X_1Y_0& X_1^2\end{array}\right).$$
However, the analogy breaks down at this stage because of the non–commutativity
$$T^{}(VV)T\left(\begin{array}{cc}\mathrm{𝟏}& \\ & \mathrm{\Phi }_1(V)\end{array}\right)$$
(259)
for (189). We leave it to the readers. There is no (well–known) direct metnod to obtain $`\mathrm{\Phi }_1(V)`$ at the current time.
Last, let us make a commemnt. For the matrix $`T`$ coming from the Clebsch–Gordan decomposition (see )
$$T=\left(\begin{array}{cccccccc}0& 0& 0& 0& 1& 0& 0& 0\\ \frac{1}{\sqrt{2}}& 0& \frac{1}{\sqrt{6}}& 0& 0& \frac{1}{\sqrt{3}}& 0& 0\\ \frac{1}{\sqrt{2}}& 0& \frac{1}{\sqrt{6}}& 0& 0& \frac{1}{\sqrt{3}}& 0& 0\\ 0& 0& 0& \frac{\sqrt{2}}{\sqrt{3}}& 0& 0& \frac{1}{\sqrt{3}}& 0\\ 0& 0& \frac{\sqrt{2}}{\sqrt{3}}& 0& 0& \frac{1}{\sqrt{3}}& 0& 0\\ 0& \frac{1}{\sqrt{2}}& 0& \frac{1}{\sqrt{6}}& 0& 0& \frac{1}{\sqrt{3}}& 0\\ 0& \frac{1}{\sqrt{2}}& 0& \frac{1}{\sqrt{6}}& 0& 0& \frac{1}{\sqrt{3}}& 0\\ 0& 0& 0& 0& 0& 0& 0& 1\end{array}\right)$$
it is not difficult to see
$`T^{}(AAA)T`$
$`=`$ $`\left(\begin{array}{cccccccc}\alpha & \overline{\beta }& & & & & & \\ \beta & \overline{\alpha }& & & & & & \\ & & \alpha & \overline{\beta }& & & & \\ & & \beta & \overline{\alpha }& & & & \\ & & & & \alpha ^3& \sqrt{3}\alpha ^2\overline{\beta }& \sqrt{3}\alpha \overline{\beta }^2& \overline{\beta }^3\\ & & & & \sqrt{3}\alpha ^2\beta & (|\alpha |^22|\beta |^2)\alpha & (2|\alpha |^2|\beta |^2)\overline{\beta }& \sqrt{3}\overline{\alpha }\overline{\beta }^2\\ & & & & \sqrt{3}\alpha \beta ^2& (2|\alpha |^2|\beta |^2)\beta & (|\alpha |^22|\beta |^2)\overline{\alpha }& \sqrt{3}\overline{\alpha }^2\overline{\beta }\\ & & & & \beta ^3& \sqrt{3}\overline{\alpha }\beta ^2& \sqrt{3}\overline{\alpha }^2\beta & \overline{\alpha }^3\end{array}\right)`$ (268)
$`=`$ $`\left(\begin{array}{ccc}\varphi _{1/2}(A)& & \\ & \varphi _{1/2}(A)& \\ & & \varphi _{3/2}(A)\end{array}\right).`$ (272)
This means a well–known decomposition
$$\frac{1}{2}\frac{1}{2}\frac{1}{2}=\left(01\right)\frac{1}{2}=\left(0\frac{1}{2}\right)\left(1\frac{1}{2}\right)=\frac{1}{2}\frac{1}{2}\frac{3}{2}.$$
E Calculation of Some Integrals
We show some integrals.
(A) Compact case :
$$<f|g>=\frac{2(2j+1)}{2\pi }_𝐂\frac{d^2z}{(1+|z|^2)^{2j+2}}f(z)\overline{g(z)}=\underset{k=0}{\overset{2j}{}}\frac{1}{{}_{2j}{}^{}C_{k}^{}}a_k\overline{b}_k$$
(273)
for $`f(z)=_{k=0}^{2j}a_kz^k`$ and $`g(z)=_{k=0}^{2j}b_kz^k`$ in $`_J`$.
This is reduced to the equation
$$\frac{2(2j+1)}{2\pi }_𝐂\frac{d^2z}{(1+|z|^2)^{2j+2}}z^k\overline{z}^l=\delta _{kl}\frac{1}{{}_{2j}{}^{}C_{k}^{}}.$$
If we use the change of variables
$$x=\sqrt{r}\text{cos}\theta ,y=\sqrt{r}\text{sin}\theta d^2z=dxdy=\frac{1}{2}drd\theta $$
then by using integration by parts
Left hand side $`=`$ $`\delta _{kl}(2j+1){\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{r^k}{(1+r)^{2j+2}}}𝑑r`$
$`=`$ $`\delta _{kl}(2j+1){\displaystyle \frac{k}{2j+1}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{r^{k1}}{(1+r)^{2j+1}}}𝑑r`$
$`=`$ $`\mathrm{}`$
$`=`$ $`\delta _{kl}(2j+1){\displaystyle \frac{k}{2j+1}}{\displaystyle \frac{k1}{2j}}\mathrm{}{\displaystyle \frac{1}{2jk+2}}{\displaystyle \frac{1}{2jk+1}}`$
$`=`$ $`\delta _{kl}{\displaystyle \frac{k!}{(2j)(2j1)\mathrm{}(2jk+1)}}`$
$`=`$ $`\delta _{kl}{\displaystyle \frac{1}{{}_{2j}{}^{}C_{k}^{}}}.`$
(B) Non–compact case :
$$<f|g>=\frac{2(2j1)}{2\pi }_Dd^2z(1|z|^2)^{2j2}f(z)\overline{g(z)}=\underset{n=0}{\overset{\mathrm{}}{}}\frac{n!}{(2j)_n}a_n\overline{b}_n$$
(274)
for $`f(z)=_{n=0}^{\mathrm{}}a_nz^n`$ and $`g(z)=_{n=0}^{\mathrm{}}b_nz^n`$ in $`H^2`$.
This is reduced to the equation
$$\frac{2(2j1)}{2\pi }_Dd^2z(1|z|^2)^{2j2}z^k\overline{z}^l=\delta _{kl}\frac{k!}{(2j)_k}.$$
Similarly in the case of (A), we obtain
Left hand side $`=`$ $`\delta _{kl}(2j1){\displaystyle _0^1}(1r)^{2j2}r^k𝑑r`$
$`=`$ $`\delta _{kl}(2j1){\displaystyle \frac{k}{2j1}}{\displaystyle _0^1}(1r)^{2j1}r^{k1}𝑑r`$
$`=`$ $`\mathrm{}`$
$`=`$ $`\delta _{kl}(2j1){\displaystyle \frac{k}{2j1}}{\displaystyle \frac{k1}{2j}}\mathrm{}{\displaystyle \frac{1}{2j+k2}}{\displaystyle \frac{1}{2j+k1}}`$
$`=`$ $`\delta _{kl}{\displaystyle \frac{k!}{(2j)(2j+1)\mathrm{}(2j+k1)}}`$
$`=`$ $`\delta _{kl}{\displaystyle \frac{k!}{(2j)_k}}`$
by using integration by parts. |
warning/0506/hep-ph0506276.html | ar5iv | text | # On leptonic decay of a heavy quarkonium with a Higgs-boson emission
## Abstract
A leptonic $`(\overline{l}l)`$ decay of a heavy quark-antiquark bound state $`T(\overline{Q}Q)`$ with a Higgs-boson $`H`$ emission is investigated. The applying of the well-known low-energy theorem to meson-Higgs coupling allows one to estimate the probability of the decay $`T(\overline{Q}Q)\overline{l}lH`$. The only a simple version of the Standard Model extension containing two-Higgs doublet is considered.
It is well-known that some extensions of the Standard Model (SM) admit the existence of new physical bound states (hadrons), composed of heavy quarks and antiquarks including the 4th generation quarks ($`Q_4`$) .
The question of the existence of 4th generation fermions ($`f_4`$) is among the most important, intriguing and not solved yet one in the modern elementary particle physics. We know that, e.g., the heterotic string phenomenology in $`E_6`$ model leads to the 4th generation of leptons ($`l_4`$) and $`Q_4`$ with a relatively stable massive neutrino of 4th generation ($`\nu _4`$) . A possible virtual contributions of 4th generation particles have been advocated by recent analysis of precision data on the SM parameters. The following question arises: what about the recent limits on the masses of $`f_4`$ ? It turns out that $`l_4`$ and $`Q_4`$ are not excluded under the condition that the Dirac $`\nu _4`$ is a (quasi)stable particle and it has a mass around 50 GeV , and the rest of a spectrum of $`f_4`$-particles satisfies their direct experimental constraints on the masses $`m_4`$ on the level above 80-220 GeV. At the moment, the best result on the lower bound restriction on $`m_4`$ was given by the CDF Collaboration at the Fermilab Tevatron, using the measurement of the energy loss $`dE/dx`$ in a ”calorimeter”. For the $`up`$-type $`Q_4`$ (labeled as $`U`$) with the electric charge $`e_U=+2/3`$ this limit is $`m_U>`$ 220 GeV , that, in principal, corresponds to the production cross-section of the order of $`1pb`$ at the Tevatron energy. It has been already reported , that in spite of the multi-$`fb^1`$ luminosity which one expects the Tevatron CDF and D0 to collect by the time the LHC will start, the rates for heavy quarks will allow their abundant production already with typical start-up luminosities of 1 $`\%`$ of the design, namely $`10^{32}cm^2s^1`$. The estimation leads to that the rate for pairs of heavy quarks production at the LHC with the mass $`O(400GeV)`$ is more than 100 times larger than at the Tevatron. In paper , we have already investigated the issues of production and decays of hadrons containing the so-called light $`Q_4`$ with the masses exceeded the top-quark mass, $`m_4>m_t`$. We considered strongly bound states, made out of heavy quarks (including fourth family) and using Higgs fields to bind them. There is an important special feature, because unlike the exchange of gauge fields, the scalar particles attract both particles and antiparticles, and the attraction of quarks by Higgs exchange is independent of color. The scenario on the hypothesis that a bound state can be formed from 6 top quarks and 6 anti-top quarks, held together mainly by Higgs particle exchange, has been considered in .
Since there is no direct indications on the existence of the stable $`f_4`$-fermions (that means their small lifetime compared with the lifetime of the Universe) it means, obviously, that one of the ways to explore these new particles is their search for via the production and their identification through the decays at modern hadron colliders. We assume that hadrons composed of $`Q_4`$-quarks are unstable and effectively decay where one of the final states should be the Higgs-boson. The reason of the Higgs-emission is covered by the more probable and effective couplings between the Higgs-boson with heavy quarks.
In this letter, we consider the process of the Higgs-boson emission in decays $`T(\overline{Q}Q)V^{}H\overline{l}lH`$, where $`T(\overline{Q}Q)`$ is the spin-1 heavy particle and $`V^{}`$ is a set of intermediate neutral vector bosons including new generations of gauge bosons (e.g., from $`E_6`$-model, Little Higgs model ). Assuming an infinitely small momentum of the Higgs-boson when the Higgs field is considered as the external one and does not carry the dependence on the coordinates (the low-energy theorem ), the probability of the decay $`T(\overline{Q}Q)\overline{l}lH`$ normalized to the Drell-Yan process $`T(\overline{Q}Q)\overline{l}l`$ is given by the formula:
$`R_{T(\overline{Q}Q)\overline{l}lH/\overline{l}l}{\displaystyle \frac{\mathrm{\Gamma }(T(\overline{Q}Q)\overline{l}lH)}{\mathrm{\Gamma }(T(\overline{Q}Q)\overline{l}l)}}={\displaystyle _0^{s_l^{max}}}𝑑s_l{\displaystyle \frac{\lambda ^{1/2}(m_T^2,m_H^2,s_l)}{24\pi ^2v^2s_l}}\eta _{HQ}^2`$ (1)
$`\left(1{\displaystyle \frac{4m_l^2}{s_l}}\right)^{1/2}\left(1+{\displaystyle \frac{2m_l^2}{s_l}}\right){\displaystyle \frac{\lambda (m_T^2,m_H^2,s_l)+6m_T^2s_l}{(m_T^2s_l)^2+\mathrm{\Gamma }_T^2m_T^2}},`$ (2)
where $`s_l=(p_l+p_{\overline{l}})^2=2(m_l^2+p_lp_{\overline{l}})`$ is the invariant mass of $`\overline{l}l`$-pair; the constants $`\eta _{HQ}`$ are defined for up (U)- and down (D)- types of quarks in the form
$`\eta _{HU}={\displaystyle \frac{\mathrm{cos}\alpha }{\mathrm{sin}\alpha }}=\mathrm{sin}(\beta \alpha )+\mathrm{cot}\beta \mathrm{cos}(\beta \alpha ),`$ (3)
$`\eta _{HD}={\displaystyle \frac{\mathrm{sin}\alpha }{\mathrm{cos}\beta }}=\mathrm{sin}(\beta \alpha )\mathrm{tan}\beta \mathrm{cos}(\beta \alpha ).`$ (4)
In the decoupling regime reflecting the special ratio between the masses of $`Z`$-boson ($`m_Z`$) and CP-odd Higgs-boson $`A`$ ($`m_A`$), $`z=(m_Z/m_A)^2<<1`$, the relations (3) and (4) transform in the following distributions on the angle $`\mathrm{tan}\beta =v_U/v_D`$ for two vacuum expectation values $`v_U`$ and $`v_D`$:
$`\eta _{HU}1+z\mathrm{sin}(2\beta )\mathrm{cos}(2\beta )\mathrm{tan}^1(\beta ),`$ (5)
$`\eta _{HD}=1z\mathrm{sin}(2\beta )\mathrm{cos}(2\beta )\mathrm{tan}(\beta ).`$ (6)
The production rate for a light CP-even Higgs-boson can be estimated. For illustration we plotted in Fig.1 in detail the $`\sqrt{s_l}`$-dependence of $`F_{T(\overline{Q}Q)\overline{l}lH/\overline{l}l}=\mathrm{\Gamma }^1(T(\overline{Q}Q)\overline{l}l)d\mathrm{\Gamma }(T(\overline{Q}Q)\overline{l}lH)/ds_l`$ on different values of $`T(Q_4\overline{Q}_4`$\- heavy quarkonia masses for $`\mathrm{tan}\beta =5`$ and $`\mathrm{tan}\beta =20`$ at fixed values of the lightest CP-even Higgs-boson mass $`m_H=m_h=`$120 GeV and $`m_A=`$300 GeV, obeying the decoupling regime above mentioned.
We found that for $`T(\overline{U}U)`$-bound state the changing of $`F_{T(\overline{Q}Q)\overline{l}lH/\overline{l}l}`$ is very small with increasing of $`m_A`$ from 200 GeV up to 300 GeV at $`5\mathrm{tan}\beta 20`$. However, the situation changes drastically if one considers the bound state composed of $`D`$-quarks. The amplitude $`F_{T(\overline{Q}Q)\overline{l}lH/\overline{l}l}`$ falls down with increasing of $`M_A`$.
In Fig. 2 we plotted the relative decay width $`R_{T(\overline{Q}Q)\overline{l}lH/\overline{l}l}`$ (1) versus the $`T`$-bound state mass $`m_T2m_Q`$ for $`T(\overline{U}U)`$\- and $`T(\overline{D}D)`$\- bound states, respectively at different $`\mathrm{tan}\beta `$.
No essential difference are found with increasing of $`\mathrm{tan}\beta `$ from 5 to 20. For conclusion, the decays of heavy quarkonia are very good place to search for a Higgs-boson in the light sector (e.g., CP-even $`h`$-boson). The decays we have discussed, $`T(\overline{Q}Q)\overline{l}lH`$, have branching ratios which could potentially be probed by precision measurements at hadron colliders. On the other hand, since there are three-body decays, the measurements of the invariant mass spectra of leptons recoiling against a Higgs-boson may give valuable insight into the dynamics of heavy quark-antiquark bound state involved. |
warning/0506/math0506513.html | ar5iv | text | # Friendly measures, homogeneous flows and singular vectors
## 1. Introduction
The theory of diophantine approximation studies how well $`𝐱=(x_1,\mathrm{},x_n)^n`$ can be approximated by $`(p_1/q_1,\mathrm{},p_n/q_n)^n`$ of a given ‘complexity’, where this complexity is usually measured by the quantity $`\mathrm{lcm}(q_1,\mathrm{},q_n)`$. Thus one is interested in minimizing the difference, in a suitable sense, between $`q𝐱`$ and a vector $`𝐩`$, where $`𝐩^n`$ and $`q`$, with a given upper bound on $`q`$. Often one finds that certain approximation problems admit a solution for almost every $`𝐱`$, while others admit a solution for almost no $`𝐱`$; one is then interested in understanding whether the typical properties remain typical when additional restrictions are placed on $`𝐱`$.
As an example, consider the notion of a singular vector, introduced by A. Khintchine in the 1920s (see \[Kh, Ca\]). Say that $`𝐱`$ is singular if for any $`\delta >0`$ there is $`T_0`$ such that for all $`TT_0`$ one can find $`𝐩^n`$ and $`q`$ with
$$q𝐱𝐩<\frac{\delta }{T^{1/n}}\mathrm{and}q<T.$$
(1.1)
a dual form, Clearly this definition is independent on the choice of the norm. Note also that by Dirichlet’s Theorem, when $`\delta >1`$ and $``$ is chosen to be the supremum norm, the system (1.1) has a nonzero integer solution for any $`T>1`$. Thus singular vectors are often referred to as those for which Dirichlet’s theorem can be infinitely improved.
Let us say that $`𝐱`$ is totally irrational if $`1,x_1,\mathrm{},x_n`$ are linearly independent over $``$. It is not hard to see that vectors which are not totally irrational are singular, and that the converse is true for $`n=1`$. However, for $`n>1`$ Khintchine \[Kh\] proved the existence of totally irrational singular vectors. On the other hand it is straightforward to verify \[Ca, Ch. V, §7\] that Lebesgue measure of the set of singular vectors is zero. In the late 1960s H. Davenport and W. Schmidt showed \[DS\] that $`𝐱^2`$ of the form $`𝐱=(t,t^2)`$ is not singular for Lebesgue-a.e. $`t`$. This was later extended to certain classes of smooth curves and higher-dimensional submanifolds of $`^n`$ by R. Baker \[Ba1, Ba2\] and M. Dodson, B. Rynne and J. Vickers \[DRV1\] respectively; see §3 for precise statement of their results.
In this paper we consider several generalizations of the notion of a singular vector. Namely, following \[Kl1, Kl2, PV\], we attach different weights to different components of $`𝐱`$ by means of the $`𝐫`$-quasinorm
$$𝐱_𝐫\stackrel{\mathrm{def}}{=}\underset{i=1,\mathrm{},d}{\mathrm{max}}|x_i|^{1/r_i},$$
where
$$𝐫=(r_1,\mathrm{},r_n)\mathrm{with}r_i>0\text{and}\underset{i=1}{\overset{n}{}}r_i=1.$$
(1.2)
Then say that $`𝐱`$ is $`𝐫`$-singular if for any $`\delta >0`$ there is $`T_0`$ such that for all $`TT_0`$ one can find $`𝐩^n`$ and $`q`$ with
$$q𝐱𝐩_𝐫<\frac{\delta }{T}\mathrm{and}q<\delta T.$$
(1.3)
Further, for $`𝐫`$ as above and an unbounded subset $`𝒯`$ of $`_1`$, say that $`𝐱`$ is $`𝐫`$-singular along $`𝒯`$ if for any $`\delta >0`$ there is $`T_0`$ such that for all $`T𝒯[T_0,\mathrm{})`$ one can find $`𝐩^n`$ and $`q`$ satisfying (1.3). We will denote the set of $`𝐫`$-singular (along $`𝒯`$) vectors by $`\mathrm{Sing}(𝐫)`$ and $`\mathrm{Sing}(𝐫,𝒯)`$ respectively.
It is clear that $`𝐱`$ is singular if and only if $`𝐱\mathrm{Sing}(𝐧)`$, where
$$𝐧\stackrel{\mathrm{def}}{=}(1/n,\mathrm{},1/n),$$
the vector assigning equal weights to each coordinate, and that $`\mathrm{Sing}(𝐫)=\mathrm{Sing}(𝐫,_1)`$ is contained in $`\mathrm{Sing}(𝐫,𝒯)`$ for any $`𝒯_1`$. However an elementary modification of the proof in \[Ca, Ch. V, §7\] shows that Lebesgue measure of $`\mathrm{Sing}(𝐫,𝒯)`$ is zero for any $`𝐫`$ as in (1.2) and any unbounded $`𝒯`$. In this paper we consider the class of friendly measures on $`^n`$, originally introduced in \[KLW\] and described in detail in §3, and prove
###### Theorem 1.1.
If $`\mu `$ is a friendly measure on $`^n`$, then for any $`𝐫`$ as in (1.2) and any unbounded $`𝒯`$, $`\mu \left(\mathrm{Sing}(𝐫,𝒯)\right)=0`$.
A special case of this theorem, with $`𝐫=𝐧`$ and $`𝒯=_1`$, was announced in \[KLW\].
The class of friendly measures includes Hausdorff measures supported on various self-similar sets such as Cantor’s ternary set, Koch snowflake, Sierpinski gasket, etc. It also includes volume measures on smooth nondegenerate submanifolds of $`^n`$. We recall that $`M^n`$ is called nondegenerate if it is parameterized by a smooth map $`𝐟`$ from an open subset $`U`$ of $`^d`$ to $`^n`$ such that for Lebesgue-a.e. $`𝐱U`$ there exists $`\mathrm{}`$ such that partial derivatives of $`𝐟`$ at $`𝐱`$ up to order $`\mathrm{}`$ span $`^n`$. If $`𝐟`$ is real analytic and $`U`$ is connected, the latter condition is equivalent to $`𝐟(U)`$ not being contained in a proper affine hyperplane of $`^n`$. Thus Theorem 1.1 significantly generalizes the aforementioned result of \[DS\] about the curve $`\{(t,t^2):t\}`$, as well as additional results obtained by several authors.
Note that it is not hard to construct a friendly measure whose support does not contain any singular vectors at all; for example, the set of badly approximable vectors supports friendly measures of arbitrarily small codimension \[KW, Ur2\]. The situation is however different for volume measures on real analytic nondegenerate manifolds. Namely, we prove
###### Theorem 1.2.
Let $`M^n`$ be a real analytic submanifold of dimension at least $`2`$ which is not contained in any proper rational affine hyperplane of $`^n`$, and let $`𝐫`$ be as in (1.2). Then there exists a totally irrational $`𝐱M\mathrm{Sing}(𝐫)`$.
Our approach to Theorem 1.1 is modelled on \[KM, KLW\]: in §2 we translate the aforementioned diophantine properties of $`𝐱^n`$ into dynamical properties of certain trajectories in the homogeneous space $`G/\mathrm{\Gamma }`$, where
$$G=\mathrm{SL}(n+1,)\text{and}\mathrm{\Gamma }=\mathrm{SL}(n+1,).$$
(1.4)
Namely, we show (Proposition 2.1) that $`𝐱^n`$ is $`𝐫`$-singular along $`𝒯`$ if and only if the corresponding trajectory leaves every compact subset of $`G/\mathrm{\Gamma }`$. To control the measure of points with divergent trajectories we employ quantitative nondivergence estimates from \[KLW\], described in detail in §4. Theorem 1.2 is proved in §5; the argument is a modification of the proof of \[We, Thm. 5.2\], and is based on ideas going back to Khintchine \[Kh\].
Acknowledgements: This research was supported by BSF grant 2000247 and NSF grant DMS-0239463. We are grateful to the Max Planck Institute for its hospitality during July 2004, and in particular to Sergiy Kolyada, who organized the activity at MPI. We are also grateful to Roger Baker for useful discussions.
## 2. Dynamical interpretation of singular vectors
Let $`G`$ and $`\mathrm{\Gamma }`$ be as in (1.4), and denote by $`\pi `$ the quotient map from $`G`$ onto $`G/\mathrm{\Gamma }`$. $`G`$ acts on $`G/\mathrm{\Gamma }`$ by left translations via the rule $`g\pi (h)=\pi (gh)`$, $`g,hG`$. Define
$$\tau (𝐱)\stackrel{\mathrm{def}}{=}\left(\begin{array}{ccccc}I_n& 𝐱& & & \\ 0& 1& & & \end{array}\right),\overline{\tau }\stackrel{\mathrm{def}}{=}\pi \tau ,$$
where $`I_n`$ stands for the $`n\times n`$ identity matrix. Then, given $`𝐫`$ as in (1.2), consider the one-parameter subgroup $`\{g_t^{(𝐫)}\}`$ of $`G`$ given by
$$g_t^{(𝐫)}\stackrel{\mathrm{def}}{=}\mathrm{diag}(e^{r_1t},\mathrm{},e^{r_nt},e^t).$$
Recall that $`G/\mathrm{\Gamma }`$ is noncompact. For an unbounded subset $`A`$ of $`_+`$ and $`xG/\mathrm{\Gamma }`$, say that a trajectory $`\{g_t^{(𝐫)}x:tA\}`$ is divergent if the map $`AG/\mathrm{\Gamma }`$, $`tg_t^{(𝐫)}x`$, is proper; that is, for any compact $`KG/\mathrm{\Gamma }`$ there exists $`t_0`$ such that $`g_t^{(𝐫)}K`$ for all $`tA[t_0,\mathrm{})`$.
It was proved in \[Da, Proposition 2.12\] that $`𝐱`$ is singular if and only if the trajectory $`\{g_t^{(𝐧)}\overline{\tau }(𝐱):t0\}`$ in $`G/\mathrm{\Gamma }`$ is divergent, and in \[Kl1, Theorem 7.4\] that $`𝐱`$ is $`𝐫`$-singular if and only if the trajectory $`\{g_t^{(𝐫)}\overline{\tau }(𝐱):t0\}`$ is divergent. We generalize this correspondence one step further:
###### Proposition 2.1.
$`𝐱`$ is $`𝐫`$-singular along $`𝒯`$ if and only if the trajectory
$$\{g_t^{(𝐫)}\overline{\tau }(𝐱):t\mathrm{log}𝒯\}G/\mathrm{\Gamma }$$
(2.1)
is divergent, where $`\mathrm{log}𝒯\stackrel{\mathrm{def}}{=}\{\mathrm{log}T:T𝒯\}`$.
To prove Proposition 2.1, we need an explicit description of compact subsets of $`G/\mathrm{\Gamma }`$. Since $`\mathrm{\Gamma }`$ is the stabilizer of $`^{n+1}`$ under the action of $`G`$ on the set of lattices in $`^{n+1}`$, $`G/\mathrm{\Gamma }`$ can be identified with $`G^{n+1}`$, that is, with the set of all unimodular lattices in $`^{n+1}`$. Fix a norm $``$ on $`^{n+1}`$, and for $`\epsilon >0`$ let
$$\begin{array}{cc}\hfill K_\epsilon & \stackrel{\mathrm{def}}{=}\pi \left(\{gG:g𝐯\epsilon 𝐯^{n+1}\{0\}\}\right);\hfill \end{array}$$
(2.2)
i.e., $`K_\epsilon `$ is the collection of all unimodular lattices in $`^{n+1}`$ which contain no nonzero vector with norm less than $`\epsilon `$. By Mahler’s compactness criterion (see e.g. \[Ra, Chapter 10\]), each $`K_\epsilon `$ is compact, and for each compact $`KG/\mathrm{\Gamma }`$ there is $`\epsilon >0`$ such that $`KK_\epsilon `$.
Now take $`𝐫`$ as in (1.2) and write
$$\overline{r}=\underset{1in}{\mathrm{min}}r_i.$$
###### Lemma 2.2.
Let $``$ be the supremum norm, let $`\epsilon `$ and $`t`$ be positive numbers with $`e^{\overline{r}t}\epsilon ,`$ and denote $`T=e^t`$. Then (1.3) with $`\delta =\epsilon ^{1/\overline{r}}`$ implies $`g_t^{(𝐫)}\overline{\tau }(𝐱)K_\epsilon `$, which in turn implies (1.3) with $`\delta =\epsilon .`$
###### Proof.
Suppose (1.3) holds with $`\delta =\epsilon ^{1/\overline{r}}`$ and with $`𝐩^n,q`$. This implies that
$$e^tq=q/T<\delta <\epsilon $$
and for $`i=1,\mathrm{},d`$,
$$e^{r_it}|qx_ip_i|<e^{r_it}\delta ^{r_i}/T^{r_i}=\epsilon ^{r_i/\overline{r}}\epsilon .$$
From this one concludes that for $`𝐯=(𝐩,q)^{n+1}\{0\}`$,
$$g_t^{(𝐫)}\tau (𝐱)𝐯=\mathrm{max}\{e^tq,e^{r_1t}|p_1+qx_1|,\mathrm{},e^{r_nt}|p_n+qx_n|\}<\epsilon ,$$
so $`g_t^{(𝐫)}\overline{\tau }(𝐱)K_\epsilon .`$ The proof of the second implication is similar and is omitted. ∎
###### Proof of Proposition 2.1.
By the preceding lemma, $`𝐱`$ is $`𝐫`$-singular along $`𝒯`$ if and only if for any $`\epsilon >0`$ there is $`T_0`$ such that $`g_t^{(𝐫)}\overline{\tau }(𝐱)K_\epsilon `$ whenever $`e^t𝒯[T_0,\mathrm{})`$. The latter, in view of Mahler’s compactness criterion, is equivalent to the fact that the trajectory (2.1) eventually leaves every compact subset of $`G/\mathrm{\Gamma }`$. ∎
As an application of this dynamical approach, we can state a condition on a measure $`\mu `$ on $`^n`$ guaranteeing that it assigns measure zero to singular vectors.
###### Proposition 2.3.
Let $`\mu `$ be a measure on $`^n`$, and suppose that for $`\mu `$-a.e. $`𝐱_0^n`$ there is a ball $`B`$ centered at $`𝐱_0`$ such that $`\delta >0`$ there exist $`\epsilon >0`$ and a sequence $`t_k\mathrm{},t_k\mathrm{log}𝒯`$, with
$$\mu \left(\{𝐱B:g_{t_k}^{(𝐫)}\overline{\tau }(𝐱)K_\epsilon \}\right)<\delta \text{ for every }k.$$
(2.3)
Then $`\mu \left(\mathrm{Sing}(𝐫,𝒯)\right)=0`$.
###### Proof.
Indeed, if we take $`\{t_k\}`$ as above and let
$$B_\epsilon \stackrel{\mathrm{def}}{=}\underset{N=1}{\overset{\mathrm{}}{}}\underset{k=N}{\overset{\mathrm{}}{}}\{𝐱B:g_{t_k}^{(𝐫)}\overline{\tau }(𝐱)K_\epsilon \},$$
then (2.3) implies that $`\mu (B_\epsilon )\delta `$. But the set $`_\epsilon B_\epsilon `$, in view of Mahler’s compactness criterion, coincides with
$$\{𝐱B:g_{t_k}^{(𝐫)}\overline{\tau }(𝐱)\text{ is divergent}\},$$
and therefore has measure zero. ∎
###### Remark 2.4.
Note that even though the definition of $`K_\epsilon `$ depends on the choice of the norm $``$ in (2.2), the assumption of Proposition 2.3 is clearly independent of this choice. Thus without loss of generality we may, and will, fix a Euclidean structure on $`^{n+1}`$ and choose $``$ to be the Euclidean norm.
## 3. The inheritance problem and friendly measures
In general, given a property which holds for a typical point in $`^n`$, it is natural to inquire for which subsets (e.g. submanifolds, self-similar ‘fractals’, etc.) a typical point on the subset also satisfies the property. The prototype for such an ‘inheritance question’ was the famous conjecture of K. Mahler from the 1930s, settled three decades later by V. Sprindžuk, which led to the theory of diophantine approximation on manifolds.
As was mentioned in the introduction, the first inheritance result related to the notion of singular vectors is due to Davenport and Schmidt: they showed that almost no points (with respect to the smooth measure class) on $`M`$ are singular, where
* $`M=\{(t,t^2):t\}^2`$ \[DS, Thm. 3\].
Later, R. Baker and Y. Bugeaud proved that almost no points on $`M`$ are singular when:
* $`M=\{(t,t^2,t^3):t\}^3`$ \[Ba1, Thm. 2\];
* $`M=\{(t,\mathrm{},t^n):t\}^n`$ \[Bu, Thm. 7\];
* $`M^2`$ is a curve with continuous third derivatives and non-vanishing curvature almost everywhere \[Ba2, Thm. 2\].
The only other paper on this topic known to us is \[DRV2\], where it was proved that almost all points on a $`C^3`$ submanifold $`M`$ of $`^n`$ are not singular if
* $`M`$ has ‘two-dimensional definite curvature almost everywhere’.
The latter condition requires the dimension of $`M`$ to be at least $`2`$ (see \[DRV1\] for more detail).
The curve in (a–c) was first studied by Mahler in connection with questions about approximation of real numbers by algebraic numbers. Note that all of the above examples are special case of nondegenerate manifolds defined in the introduction (see \[KM, Remark 6.3\] for a discussion of the relation between nondegeneracy of $`M`$ and the conditions of \[DRV1, DRV2\].)
A more general framework for discussing the inheritance problem is to recast it in terms of measures. That is, given a property $`𝒫`$ which holds for Lebesgue-a.e. $`𝐱^n`$, one wants to describe measures $`\mu `$ such that $`𝒫`$ also holds for $`\mu `$-a.e. $`𝐱`$. Let us recall certain properties of a measure on $`^n`$, introduced in \[KLW\].
Suppose $`\mu `$ is a locally finite Borel measure on $`^n`$. Let $`B(𝐱,r)`$ denote the open ball of radius $`r`$ centered at $`𝐱`$. Suppose $`U^n`$ is open. We say that $`\mu `$ is $`D`$-Federer on $`U`$ if for all $`𝐱\mathrm{supp}\mu U`$ one has
$$\frac{\mu \left(B(𝐱,3r)\right)}{\mu \left(B(𝐱,r)\right)}<D$$
whenever $`B(𝐱,3r)U`$.
Say that $`\mu `$ is nonplanar if $`\mu ()=0`$ for any affine hyperplane $``$ of $`^n`$. For an affine subspace $`^n`$ we denote by $`d_{}(𝐱)`$ the (Euclidean) distance from $`𝐱`$ to $``$, and let
$$^{(\epsilon )}\stackrel{\mathrm{def}}{=}\{𝐱^n:d_{}(𝐱)<\epsilon \}.$$
If $`B^n`$ with $`\mu (B)>0`$ and $`f`$ is a real-valued function on $`^n`$, let
$$f_{\mu ,B}\stackrel{\mathrm{def}}{=}\underset{xB\mathrm{supp}\mu }{sup}|f(x)|.$$
Given $`C`$, $`\alpha >0`$ and an open subset $`U`$ of $`^n`$, say that $`\mu `$ is $`(C,\alpha )`$-decaying on $`U`$ if for any non-empty open ball $`BU`$ centered in $`\mathrm{supp}\mu `$, any affine hyperplane $`^n`$, and any $`\epsilon >0`$ one has
$$\mu \left(B^{(\epsilon )}\right)C\left(\frac{\epsilon }{d_{}_{\mu ,B}}\right)^\alpha \mu (B).$$
Finally, let us say that $`\mu `$ is friendly if it is nonplanar, and for $`\mu `$-a.e. $`x^n`$ there exist a neighborhood $`U`$ of $`x`$ and positive $`C,\alpha ,D>0`$ such that $`\mu `$ is $`D`$-Federer and $`(C,\alpha )`$ decaying on $`U`$.
The class of friendly measures is rather large; its properties are discussed in \[KLW\] and examples are given in \[KLW, KW, Ur1, SU\]. For example, it is essentially proved in \[KM\] (see \[KLW, Lemma 7.1 and Propositions 7.2, 7.3\]) that the natural measure on a nondegenerate manifold obtained by pushing forward Lebesgue measure on $`^n`$ is friendly. Thus our main result (Theorem 1.1) supersedes entries (a–e) in the above list. Further, there are many more possibilities of interesting choices of sets which can support friendly measures. A particularly nice choice is given by limit sets of finite irreducible systems of contracting similarities (or, more generally, self-conformal maps of $`^n`$) satisfying the open set condition, see \[KLW, §8\] and \[Ur1\] for more detail.
We can now state the main measure estimate from which Theorem 1.1 will easily follow.
###### Theorem 3.1.
Let $`\mu `$ be a friendly measure on $`^n`$. Then for $`\mu `$-a.e. $`𝐳^n`$ there is a ball $`B`$ centered at $`𝐳`$ and positive $`\stackrel{~}{C},\alpha `$ with the following property: for any $`𝐫`$ as in (1.2) there exists $`t_0>0`$ such that for all $`t>t_0`$ and all $`\epsilon >0`$ one has
$$\mu \left(\{𝐱B:g_t^{(𝐫)}\overline{\tau }(𝐱)K_\epsilon \}\right)<\stackrel{~}{C}\epsilon ^\alpha \mu (B).$$
(3.1)
###### Proof of Theorem 1.1 assuming Theorem 3.1.
Take $`B`$, $`\stackrel{~}{C}`$, $`\alpha `$ as in Theorem 3.1, given $`\delta >0`$ choose $`\epsilon >0`$ with $`\stackrel{~}{C}\epsilon ^\alpha \mu (B)<\delta `$, and let $`\{t_k\}`$ be any unbounded subsequence of $`\mathrm{log}𝒯(t_0,\mathrm{})`$. Then (2.3) becomes an immediate consequence of (3.1), and an application of Proposition 2.3 finishes the proof. ∎
## 4. A quantitative nondivergence estimate
We will derive Theorem 3.1 from a quantitative nondivergence result. To state it we need some additional definitions.
Let $`f:^n`$. Given $`C`$, $`\alpha >0`$, $`U^n`$ and a measure $`\mu `$ on $`^n`$, say that $`f`$ is $`(C,\alpha )`$-good on $`U`$ with respect to $`\mu `$ if for any ball $`BU`$ centered in $`\mathrm{supp}\mu `$ and any $`\epsilon >0`$ one has
$$\mu \left(\{yB:|f(y)|<\epsilon \}\right)C\left(\frac{\epsilon }{f_{\mu ,B}}\right)^\alpha \mu (B).$$
We refer the reader to \[KM, KLW\] for various properties and examples. We are going to need two elementary observations, which we state below for convenience.
###### Lemma 4.1.
Let $`U^n`$ be open, $`C,\alpha >0`$, $`\mu `$ a measure on $`^n`$.
* \[KLW, Lemma 4.2\] $`\mu `$ is $`(C,\alpha )`$-decaying on $`U`$ if and only if any affine function (equivalently, any function of the form $`d_{}`$, where $``$ is an affine hyperplane) is $`(C,\alpha )`$-good on $`U`$.
* \[KLW, Lemma 4.1\] If $`f_1,\mathrm{},f_k`$ are $`(C,\alpha )`$-good on $`U`$ w. r. t. $`\mu `$, then the function $`𝐱𝐟(𝐱)`$, where $`𝐟=(f_1,\mathrm{},f_k)`$ and $``$ is the Euclidean norm on $`^k`$, is $`(k^{\alpha /2}C,\alpha )`$-good on $`U`$ w. r. t. $`\mu `$.
Let
$$𝒲\stackrel{\mathrm{def}}{=}\text{ the set of nonzero rational subspaces of }^{n+1}.$$
For $`V𝒲`$ and $`gG`$, let
$$\mathrm{}_V(g)\stackrel{\mathrm{def}}{=}g(𝐯_1\mathrm{}𝐯_k),$$
where $`\{𝐯_1,\mathrm{},𝐯_k\}`$ is a generating set for $`^{n+1}V`$ and $``$ is the extension of the Euclidean norm from $`^{n+1}`$ to its exterior algebra; note that $`\mathrm{}_V(g)`$ does not depend on the choice of $`\{𝐯_i\}`$.
We will use the following estimate, which is a special case of \[KLW, Theorem 4.3\]:
###### Theorem 4.2.
Given $`n`$ and positive constants $`C,D,\alpha `$, there exists $`\stackrel{~}{C}=\stackrel{~}{C}(n,C,D,\alpha )>0`$ with the following property. Suppose $`\mu `$ is $`D`$-Federer on an open subset $`U`$ of $`^n`$, $`h`$ is a continuous map $`UG`$, $`0<\rho 1`$, $`𝐳U\mathrm{supp}\mu `$, and $`B=B(𝐳,r)`$ is a ball such that denote $`B(𝐱,cr)`$ $`B(𝐳,3^nr)U`$, and that for each $`V𝒲`$,
* the function $`\mathrm{}_Vh`$ is $`(C,\alpha )`$-good on $`B(𝐳,3^nr)`$ with respect to $`\mu `$,
and
* $`\mathrm{}_Vh_{\mu ,B}\rho `$.
Then for any $`0<\epsilon \rho `$,
$$\mu \left(\{𝐱B:\pi \left(h(𝐱)\right)K_\epsilon \}\right)\stackrel{~}{C}(\epsilon /\rho )^\alpha \mu (B).$$
(4.1)
###### Proof of Theorem 3.1.
Recall that we are given a friendly measure $`\mu `$. For $`\mu `$-almost every $`𝐳^n`$, choose a neighborhood $`U`$ of $`𝐳`$, positive constants $`C^{},D,\alpha `$ such that $`\mu `$ is $`D`$-Federer and $`(C^{},\alpha )`$-decaying on $`U`$, and a ball $`B=B(𝐳,r)`$ centered at $`𝐳`$ such that $`B(𝐳,3^nr)`$ is contained in $`U`$. Clearly the desired estimate (3.1) will coincide with (4.1) if one takes $`\rho =1`$ and lets $`h=h_{𝐫,t}`$ where the latter is defined by
$$h_{𝐫,t}(𝐱)\stackrel{\mathrm{def}}{=}g_t^{(𝐫)}\tau (𝐱).$$
Therefore it suffices to verify the assumptions of Theorem 4.2 for the above choice of $`\rho `$ and $`h`$. This is done below in Lemmas 4.3 and 4.4. ∎
###### Lemma 4.3.
Suppose that $`\mu `$ is $`(C^{},\alpha )`$-decaying on an open $`U^n`$. Then for any $`V𝒲`$, any $`𝐫`$ as in (1.2) and any $`t0`$,
$$\text{the function }\mathrm{}_Vh_{𝐫,t}\text{ is }(C,\alpha )\text{-good on }U\text{ with respect to }\mu ,$$
where $`C=(n+1)^{\alpha /2}C^{}`$.
###### Lemma 4.4.
Suppose that $`\mu `$ is nonplanar. Then for any $`𝐫`$ as in (1.2) and any ball $`B`$ with $`\mu (B)>0`$ there is $`t_0=t_0(\mu ,𝐫,B)`$ such that for any $`V𝒲`$ and any $`tt_0`$ one has
$$\mathrm{}_Vh_{𝐫,t}_{\mu ,B}1.$$
The proof of both lemmas hinges on a computation of the $`h_{𝐫,t}(𝐱)`$-action on the exterior powers of $`^{n+1}`$, as in \[KM\] or \[KLW\]. We include the argument for the sake of completeness. For the remainder of the section, to simplify notation we will write $`g_t`$ instead of $`g_t^{(𝐫)}`$ and $`h_t`$ instead of $`h_{𝐫,t}`$.
Denote by $`V_0`$ the subspace
$$V_0\stackrel{\mathrm{def}}{=}\{(x_1,\mathrm{},x_{n+1}):x_{n+1}=0\}$$
of $`^{n+1}`$, and let $`𝐞_0\stackrel{\mathrm{def}}{=}(0,\mathrm{},0,1)`$ be a vector orthonormal to $`V_0`$. Note that for any $`𝐱^n`$,
$$\tau (𝐱)\text{ acts trivially on }V_0\text{ and }\tau (𝐱)𝐞_0=𝐞_0+𝐱,$$
(4.2)
where we with some abuse of notation identified $`V_0`$ with $`^n`$.
Now suppose that $`V`$ is a $`k`$-dimensional subspace of $`^{n+1}`$, $`k1`$, spanned by integer vectors $`𝐯_1,\mathrm{},𝐯_k`$, and denote $`𝐯_1\mathrm{}𝐯_k`$ by $`𝐰`$. By applying Gaussian elimination over the integers to $`\{𝐯_1,\mathrm{},𝐯_k\}`$ one can write $`𝐰`$ in the form
$$𝐰=𝐰_0(q𝐞_0𝐩),$$
(4.3)
where $`q`$, $`𝐩V_0()`$ and $`𝐰_0^{k1}\left(V_0()\right)`$. Using (4.2) and (4.3), one writes
$$\tau (𝐱)𝐰=𝐰_0\left(q𝐞_0+q𝐱𝐩\right)=𝐰_0(q𝐱𝐩)+q𝐰_0𝐞_0,$$
and hence
$$h_t(𝐱)𝐰=g_t\left(𝐰_0(q𝐱𝐩)\right)+qg_t(𝐰_0𝐞_0).$$
(4.4)
Note that the two summands in (4.4) are orthogonal, therefore
$`\left(\mathrm{}_Vh_t(𝐱)\right)^2`$ $`=h_t(𝐱)𝐰^2`$ (4.5)
$`=q^2g_t(𝐰_0𝐞_0)^2+g_t\left(𝐰_0(q𝐱𝐩)\right)^2.`$
Now we can return to the lemmas.
###### Proof of Lemma 4.3.
Write the second summand in (4.5) in the form
$$g_t𝐰_0g_t(q𝐱𝐩)^2=g_t𝐰_0^2d_{g_t𝒫}\left(g_t(q𝐱𝐩)\right)^2,$$
where $`𝒫`$ stands for the linear subspace of $`V_0`$ corresponding to $`𝐰_0`$. For any $`t`$, $`q`$ and $`𝐩`$, the function $`𝐱d_{g_t𝒫}^2\left(g_t(q𝐱𝐩)\right)`$ is the sum of squares of at most $`n`$ affine functions, each of which is $`(C^{},\alpha )`$-good on $`U`$ with respect to $`\mu `$ in view of Lemma 4.1(a) and the assumption of Lemma 4.3. Therefore, by Lemma 4.1(b), the function $`\mathrm{}_Vh_t`$ is $`((n+1)^{\alpha /2}C^{},\alpha )`$-good on $`U`$ with respect to $`\mu `$. ∎
###### Proof of Lemma 4.4.
Let us denote by $`e^{\gamma t}`$ the smallest eigenvalue of the induced action of $`g_t`$ on $`^k(V_0)`$ (here $`\gamma >0`$ depends on $`𝐫`$). If $`q=0`$, in view of (4.3) we have $`𝐰^k\left(V_0()\right)`$, hence
$$\mathrm{}_Vh_t(𝐱)=g_t\tau (𝐱)𝐰=g_t𝐰e^{\gamma t}𝐰1$$
for all $`t0`$. Thus the conclusion of the lemma holds with e.g. $`t_0=0`$. Otherwise, using (4.4), one can write
$`\mathrm{}_Vh_t(𝐱)`$ $`g_t\left(𝐰_0(q𝐱𝐩)\right)e^{\gamma t}𝐰_0(q𝐱𝐩)`$
$`=e^{\gamma t}𝐰_0d_𝒫(q𝐱𝐩)=|q|e^{\gamma t}𝐰_0d_{(𝒫+𝐩/q)}(𝐱)e^{\gamma t}d_{(𝒫+𝐩/q)}(𝐱)`$
(the last inequality holds since both $`q`$ and all the coordinates of $`𝐰_0`$ are integers).
If $`B`$ is a ball with $`\mu (B)>0`$, an easy compactness argument using the assumption that $`\mu `$ is nonplanar shows the existence of $`c=c(B)>0`$ such that $`d_{}_{\mu ,B}c`$ for any proper affine subspace $``$ of $`^n`$. Hence
$$\mathrm{}_Vh_t_{\mu ,B}c|q|e^{\gamma t}𝐰_0ce^{\gamma t},$$
and the conclusion of the lemma holds with $`t_0=\frac{1}{\gamma }\mathrm{log}\frac{1}{c}`$. ∎
## 5. Constructing singular vectors on submanifolds
In this section we adapt the methods of the paper \[We\] and exhibit points with divergent trajectories on certain proper subsets of $`G/\mathrm{\Gamma }`$.
###### Theorem 5.1.
Let $`\{g_t:t\}`$ be a one-parameter subgroup of $`G`$. Suppose $`X`$ is a closed subset of $`G`$, and $`\{X_i:i\}`$ and $`\{X_j^{}:j\}`$ are two lists of subsets of $`X`$, such that for some strictly decreasing continuous $`\psi :_+_+`$, the following conditions are satisfied.
1. Density. For every $`j`$, $`X_j=\overline{X_j_{ij}X_i}.`$
2. Transversality I. For every $`ij`$, $`X_i=\overline{X_iX_j}.`$
3. Transversality II. For any $`i,j`$, $`X_i=\overline{X_iX_j^{}}.`$
4. Local Uniformity w.r.t. $`\{K_{\psi (t)}\}`$. For every $`i`$ and every $`xX_i`$ there is a neighborhood $`U`$ in $`G`$ and $`t_0`$ such that for all $`tt_0`$ and all $`zUX_i`$,
$$g_t\pi (z)K_{\psi (t)}.$$
Then there is $`xX\left(_iX_i_jX_j^{}\right)`$ and $`t_0>0`$ such that for all $`tt_0`$, $`g_t\pi (x)K_{\psi (t)}.`$
###### Proof.
Equip $`X`$ with the relative topology inherited from $`G`$, and write
$$K(t)=K_{\psi (t)}.$$
We will construct inductively a sequence of open sets with compact closure $`\mathrm{\Omega }_0,\mathrm{\Omega }_1,\mathrm{\Omega }_2,\mathrm{}`$ in $`X`$, an increasing sequence of indices $`i_1,i_2,\mathrm{}`$, and an increasing sequence of times $`T_0,T_1,\mathrm{},`$ such that the following hold for $`k=1,2,\mathrm{}`$:
* $`\overline{\mathrm{\Omega }_k}\mathrm{\Omega }_{k1}`$.
* For every $`j<i_k`$, $`X_j\mathrm{\Omega }_k=\mathrm{}`$ and $`X_j^{}\mathrm{\Omega }_k=\mathrm{}`$.
* $`X_{i_k}\mathrm{\Omega }_k`$ is nonempty and for every $`zX_{i_k}\mathrm{\Omega }_k`$ and every $`tT_k`$ we have $`g_t\pi (z)K(t).`$
We will also have for $`k=2,3,\mathrm{}`$:
* For every $`z\mathrm{\Omega }_k`$ and every $`t[T_{k1},T_k]`$, $`g_t\pi (z)K(t)`$.
To see that such sequences suffice, note that by condition a, $`_k\mathrm{\Omega }_k`$ is nonempty, and for $`z\mathrm{\Omega }_k`$ we have by condition b that $`z_iX_i_jX_j^{}`$ and by condition d that $`g_t\pi (z)K(t)`$ for $`tT_1`$.
Now let us construct the sequences inductively. Choose $`T_0=0`$, $`i_1=1`$. Let $`xX_1`$ and using the local uniformity hypothesis, let $`\mathrm{\Omega }_1`$ be a small enough open neighborhood of $`x`$, and $`T_1`$ large enough, so that for all $`zX_1\mathrm{\Omega }_1`$ and all $`tT_1,`$ we have $`g_t\pi (z)K(t).`$ Now letting $`T_0`$ be arbitrary and $`\mathrm{\Omega }_0`$ be any open set containing $`\overline{\mathrm{\Omega }_1}`$, we see that a, b and c hold for $`k=1`$.
Suppose we have chosen $`i_s,\mathrm{\Omega }_s,T_s`$ for $`s=1,\mathrm{},k`$. By the density condition there are $`\mathrm{}i_k`$ such that
$$X_{\mathrm{}}\mathrm{\Omega }_kX_{i_k}\mathrm{}.$$
Choose for $`i_{k+1}`$ any such $`\mathrm{}`$. Note that $`i_{k+1}>i_k`$ by b. Let $`xX_{i_k}\mathrm{\Omega }_kX_{i_{k+1}}.`$ By the local uniformity assumption, there is a small enough open neighborhood $`U`$ of $`x`$ and a large enough $`T_{k+1}`$ such that for all $`zUX_{i_{k+1}}`$ and all $`tT_{k+1}`$, $`g_t\pi (z)K(t).`$ In addition let $`U`$ be small enough so that $`\overline{U}\mathrm{\Omega }_k`$. Since $`xX_{i_k}`$, $`g_t\pi (x)K(t)`$ for $`t[T_k,T_{k+1}]`$, hence by continuity of $`\psi `$ and the action, there is a small enough neighborhood $`\stackrel{~}{\mathrm{\Omega }}`$ of $`x`$ contained in $`U`$ so that
$$z\stackrel{~}{\mathrm{\Omega }},t[T_{k1},T_k]g_t\pi (z)K(t).$$
Now we can define $`\mathrm{\Omega }_{k+1}`$ by
$$\mathrm{\Omega }_{k+1}=\stackrel{~}{\mathrm{\Omega }}\underset{j<i_{k+1}}{}\left(X_jX_j^{}\right).$$
We now verify that $`i_{k+1},\mathrm{\Omega }_{k+1},T_{k+1}`$ satisfy the required conditions. Condition a holds by our choice of $`U.`$ Condition b follows from the definition of $`\mathrm{\Omega }_{k+1}`$. In condition c, $`\mathrm{\Omega }_{k+1}X_{i_{k+1}}\mathrm{}`$ because $`x\stackrel{~}{\mathrm{\Omega }}X_{i_{k+1}}`$, and because of the transversality assumptions. The second assertion in condition c holds because of the choice of $`T_{k+1}`$ and $`U`$. Condition d holds because of the choice of $`\stackrel{~}{\mathrm{\Omega }}`$. ∎
We now derive a consequence of this theorem. This requires some notation. Fix $`𝐫`$, let $`M`$ be a submanifold of $`^n`$, and choose $`1k<\mathrm{}n`$. For $`𝐯^n`$ and $`s`$ let
$$L_𝐯(s)=\{𝐱M:𝐱,𝐯=s\}.$$
Let $`𝐞_k,𝐞_{\mathrm{}}`$ be the $`k`$-th and $`\mathrm{}`$-th standard basis vectors and let $`\{X_i\}`$ be a list of the distinct connected components of the sets $`\{L_{𝐞_r}(s):r\{k,\mathrm{}\},s\}`$. Also let $`\{X_j^{}\}`$ be a list of the distinct connected components of the sets $`\{L_𝐯(s):𝐯^n,s\}`$ which are not in the list $`\{X_i\}`$. In case $`\{X_j^{}\}=\{X_1^{},\mathrm{},X_m^{}\}`$ happens to be a finite list, we put $`X_j^{}=\mathrm{}`$ for $`j>m`$, i.e., we may assume that $`\{X_j^{}\}`$ is indexed by $``$ as well.
###### Corollary 5.2.
Suppose $`𝐫,M,k,\mathrm{},\{X_i\},\{X_j^{}\}`$ are as above, and that the density and two transversality assumptions of Theorem 5.1 are satisfied. Suppose also that $`\delta :_+_+`$ is a function such that
$$T^\rho \delta (T)_T\mathrm{}\mathrm{},\mathrm{where}\rho =\mathrm{min}\{r_k,r_{\mathrm{}}\}/n.$$
Then there is a totally irrational $`𝐱M`$ and $`T_0`$ such that for all $`TT_0`$ there is $`𝐩^n,q`$ satisfying (1.3) for $`\delta =\delta (T).`$
###### Proof.
Set $`\psi (t)\stackrel{\mathrm{def}}{=}\delta (e^t)`$, so that
$$e^{\rho t}\psi (t)\mathrm{}.$$
(5.1)
Suppose without loss of generality that $`\psi `$ is a decreasing function. Using Lemma 2.2, it suffices to show that there are $`x=\tau (𝐱)\tau (M)`$ and $`t_0`$ such that $`𝐱`$ is totally irrational and for all $`tt_0`$,
$$g_t^{(𝐫)}\pi (x)K_{\psi (t)}.$$
To this end, we apply Theorem 5.1 to the lists $`\{X_i\},\{X_j^{}\}`$ (which we identify with their images under $`\tau `$, and thus consider them as subsets of $`G`$). Since we have assumed the density and two transversality conditions, we need only check the locally uniform escape condition.
To verify the condition of local uniformity w.r.t. $`\{K_{\psi (t)}\}`$, fix $`X_i`$ so that (possibly after exchanging $`k`$ and $`\mathrm{}`$), for some $`p/q`$ and for all $`𝐳=(z_1,\mathrm{},z_n)X_i`$ we have $`z_k=p/q`$. Let
$$W=\{(w_1,\mathrm{},w_{n+1})^{n+1}:qw_k+pw_{n+1}=0\}.$$
This is a rational linear subspace of $`^{n+1}`$ of dimension $`n`$, with the property that for all $`𝐳X_i`$, all vectors in $`\tau (𝐳)W`$ have their $`k`$-th coordinate equal to zero. It follows that for any bounded neighborhood $`U`$ intersecting $`X_i`$ there is a constant $`C`$, such that for all $`𝐳UX_i`$ we have
$$\mathrm{}_W\left(g_t^{(𝐫)}\tau (𝐳)\right)Ce^{r_kt}.$$
Using Minkowski’s convex body theorem, we find a constant $`C^{}`$ so that for all $`t`$ and all $`𝐳UX_i`$, $`g_t^{(𝐫)}\tau (𝐳)\left(W^{n+1}\right)`$ contains a non-zero vector of length at most
$$C^{}e^{r_kt/n}C^{}e^{\rho t}.$$
Now from (5.1) if follows that there is $`t_0`$ such that for all $`tt_0`$, the length of such a vector is less than $`\psi (t)`$. This concludes the proof. ∎
###### Proof of Theorem 1.2.
For $`𝐳M`$ let $`T_𝐳M^n`$ be the tangent space to $`M`$ at $`𝐳`$. Since $`dim(M)2`$, there are distinct indices $`k,\mathrm{}\{1,\mathrm{},n\}`$ and an open $`VM`$ such that if $`P`$ is the projection
$$P(𝐱)=(x_k,x_{\mathrm{}})\mathrm{for}𝐱=(x_1,\mathrm{},x_n),$$
then the derivative $`D_𝐳\left(P|_M\right):T_𝐳M^2`$ is surjective for all $`𝐳V`$. With no loss of generality, replace $`M`$ with $`V`$ and define $`\{X_i\},\{X_j^{}\}`$ as in the paragraph preceding Corollary 5.2. In view of the Corollary, it remain to check that the density and two transversality hypotheses hold for $`\{X_i\},\{X_j^{}\}`$.
Since these hypotheses hold for horizontal and vertical lines in an open subset of the plane, and since the $`X_i`$’s are the pre-images of these lines under $`P`$, we see that the density and the first transversality hypothesis hold. Now for $`i`$ and $`j`$, suppose that $`X_iX_j^{}\mathrm{}`$ (otherwise there is nothing to prove). Then $`X_i`$ and $`X_iX_j^{}`$ are connected analytic submanifolds of $`M`$. Suppose if possible that $`X_iX_j^{}`$ contains an open subset of $`X_i`$. Then, since they are analytic and $`X_i`$ is connected, $`X_iX_j^{}`$. Since $`M`$ is not contained in a rational affine hyperplane, both $`X_i`$ and $`X_j^{}`$ are submanifolds of $`M`$ of codimension one, and since $`X_j^{}`$ is also connected we must have $`X_i=X_j^{}`$, contrary to the construction. This implies that $`X_i=\overline{X_iX_j^{}}`$, as required. ∎
###### Remark 5.3.
The hypotheses of Theorem 1.2 are clearly satisfied when $`M`$ is an analytic nondegenerate submanifold of dimension at least $`2`$, and also when $`M`$ is an affine subspace of dimension at least $`2`$ not contained in any rational affine hyperplane.
###### Remark 5.4.
The proof of Theorem 1.2 actually shows that $`M\mathrm{Sing}(𝐫)`$ contains uncountably many totally irrational vectors. Indeed, given any countable subset $`A=\{𝐳_1,𝐳_2,\mathrm{}\}`$ of $`M`$, replace the list $`\{X_j^{}\}`$ with the list $`\{X_1^{},\{𝐳_1\},X_2^{},\{𝐳_2\},\mathrm{}\}`$. Applying the same argument yields an element of $`M\mathrm{Sing}(𝐫)`$ which is totally irrational and does not belong to $`A`$. |
warning/0506/cond-mat0506635.html | ar5iv | text | # Mean-field limit of systems with multiplicative noise
## I Introduction
Problems susceptible to be mathematically represented by stochastic (Langevin) equations including a multiplicative noise abound not only in physics, but also in biology, ecology, economy, or social sciences. In a broad sense a Langevin equation is said to be multiplicative if the noise amplitude depends on the state variable/s itself/themselves vK . In this sense, problems exhibiting absorbing states, i.e. fluctuation-less states in which the system can be trapped, are described by equations whose noise amplitude is proportional to the square-root of the (space and time dependent) activity density, vanishing at the absorbing state AS . Systems within this class are countless: propagating epidemics, autocatalytic reactions, reaction-diffusion problems, self-organized criticality, pinning of flux lines in superconductors, etc. AS .
In a more restrictive sense, the one we will use from now on, it is customary to restrict the term multiplicative noise (MN) to noise amplitudes linear in the activity density Graham ; Sancho ; Redner ; Sornette . Such equations appear ubiquitously in economics, optics, population dynamics, study of instabilities, etc. In all these cases, multiplicative noise terms appear rather straightforwardly when constructing (more or less rigorously) stochastic representative equations. In their spatially extended version, Langevin equations including a MN (see Eq. (1) below) were first proposed, to the best of our knowledge, in the context of synchronization of coupled map lattices PK , and soon after studied in Refs. MN ; Walter ; vdB ; Marsili ; Birner . Such equations describe different situations; in particular, there has been a recent interest in their application to non-equilibrium wetting Lisboa and also to synchronization problems in extended systems Synchro ; see Ref. MNreview for a recent review. In all these cases, there is a phase transition between an active phase in which the system has a stationary non-vanishing activity and another one in which the density field falls continuously toward zero without ever reaching it in finite time. The analogies and differences between this family of models and the one with a square-root type of noise, in which the absorbing state is reachable within a finite time, have been discussed in nature ; Janssen .
It is remarkable that by using the so-called Cole-Hopf transformation, the spatially extended MN Langevin equation can be mapped into a non-equilibrium interface, as represented by the Kardar-Parisi-Zhang (KPZ) equation KPZ ; HZ in the presence of a bounding wall MNreview . In this language, the active phase describes interfaces pinned by the wall while the absorbing one corresponds to depinned KPZ-like interfaces moving away from the wall.
In recent years many aspects of systems with MN have been elucidated. For example, a renormalization group approach has been constructed, scaling relations derived, and the large-$`N`$ limit (where $`N`$ is the number of components) studied MN . It is well established that above two-dimensions there are two different regimes; a weak-noise and a strong-noise one, in full analogy with the known phenomenology of KPZ interfaces HZ . There is however a crucial point which remains to be fully understood: the mean-field behavior of such systems. Another interesting mapping is that in the absence of the non-linear term, the MN equation corresponds to the equation governing the evolution of directed polymers in random media (See HZ and references therein, as well as MNreview ).
In Langevin equations with additive noise, as for instance the Model A describing the universality class of Ising like transitions HH , the mean-field approximation can be obtained in a number of equivalent ways, all of them leading to the same results with different degrees of sophistication. For example, mean-field critical exponents can be obtained (i) by discarding the noise and solving the remaining deterministic equation, (ii) as the most probable solution in a path integral formulation, (iii) self-consistently by assuming that each site sees the average of the remaining, (iv) by naive power counting in the corresponding action, (v) as the lowest order in a perturbative loop expansion, etc.
Contrarily, in systems with MN the meaning of the mean-field solution is much less clear. Indeed, early studies Walter led to conflicting results depending on the considered approximation method. For example, by removing the noise, much of the physics is lost and trivial results (identical to those for additive noise) are obtained. Also, if all spatial dependence is eliminated (by removing the Laplacian term in the MN equation), one is left with a solvable zero-dimensional equation which does not reproduce faithfully the rich mean-field phenomenology. Therefore, contrarily to more standard problems, both space-dependence and noise have to be retained in order to construct a sound mean-field solution.
If the bounding wall is eliminated then one is left with the directed polymer in random media equation, for which many results are available. In particular, Derrida and collaborators worked out a solution on the Cayley tree, mean-field results, $`1/d`$ expansions, and solutions on hierarchical lattices exits Derrida . Also, Mézard and Parisi derived a variational approach in replica space MP , and also Fisher and collaborators reached also interesting results on these issues Fisher . For a more detailed review on results for directed-polymers see HZ and references therein.
On the other hand, restoring the bounding wall (which is the case we are interested in) Birner et al. Birner (see also Walter ; vdB ; Marsili ) performed a self-consistent calculation and reported on the existence of two different mean-field behaviors: a weak-noise regime (in this particular case, the noise can be completely disregarded) and a strong-noise one exhibiting non-universal exponents depending on the noise amplitude. This is in agreement with the field theoretic expectation of two different behaviors for MN-like and KPZ-like equations in high dimensional systems (where mean-field results are expected to be valid) HZ ; MNreview . These results are difficult to compare with the abovementioned ones obtained for directed polymers in random media, which correspond to the deppined phase of the full problem, while the self-consistent approach is intrinsically devised for the pinned phase.
Let us also underline that a full understanding in terms of path integrals and extremal paths is still missing despite some efforts in this direction Kharchenko ; Fogedby .
Together with the determination of the right mean-field theory, another relevant and highly debated issue is to establish the upper critical dimension $`d_c`$, above which mean-field results hold. As some exponents, as the dynamical one $`z`$ have been claimed MN ; MNreview to coincide for MN and KPZ, both problems are expected to have the same upper critical dimension. While there is consensus that the mean-field weak-noise regime should be valid above $`d_c=2`$ (coinciding with the critical dimension for weak-noise Edwards-Wilkinson interfaces HZ ) there exist highly conflicting results (some pointing to $`d_c=4`$ KPZdc4 and some supporting $`d_c=\mathrm{}`$ KPZdcinf ) for the strong-noise one regime.
In any generic problem, for sufficiently high dimensions every site in a spatially extended system “sees” the average of its neighbors (assuming the space has been discretized) which can be correctly approximated by its mean value (the distribution of the average values is well described by the standard central limit theorem above the upper critical dimension) defining in this way a sound mean-field solution. In the case of MN, as we will show, the situation is somehow anomalous: in the strong-noise regime the mean-field solution itself breaks down in the neighborhood of the critical point once fluctuations are taken into account (as a version of the Ginzburg criterion shows). This stems from the primary role played by fluctuations and rare events in multiplicative processes Redner ; Sornette , and makes one wonder whether a mean-field solution can be valid at all in any finite dimension.
In this paper we revisit the self-consistent mean-field solution of Langevin equations with MN and report on a new previously overlooked regime. After that, we discuss under which circumstances such a solution breaks down (Ginzburg criterion Birgenau ; LB ). Also, we present a fluctuating-solution aimed to extend the self-consistent one by allowing for fluctuations of the average field. Finally, we present some speculations on the issue of the existence of an upper critical dimension for this type of systems, i.e. if there is a finite space dimensionality above which the non-fluctuating solution holds or not.
The paper is organized as follows. In Section II, we outline a self-consistent solution of the MN equation defined in an infinite fully connected lattice. Different regimes are found, corresponding to weak-noise, intermediate-noise, and strong-noise respectively. By constructing a Ginzburg criterion we will show how the strong-noise solution is intrinsically unstable as soon as fluctuations are taken into account. In Section III we study the solution in finite lattices (fluctuating solution) and show by means of numerical simulations how the mean-field solution breaks down in the vicinity of the critical point. Section IV is devoted to the study of multiplicative noise on random regular graphs. This is done in order to ascertain the type of behavior in the thermodynamic limit when the connectivity of each site remains finite and, therefore, in order to obtain some insight on the behavior of the MN equation in large but finite space dimensions. The final section contains a summary of the findings along with some concluding remarks. Finally, a general scaling theory is presented in an appendix.
## II Non Fluctuating Solution
We summarize and extend the mean-field solution obtained by Birner et al. Birner for the MN equation as described in MN ; MNreview :
$$\dot{\varphi }_i=a\varphi _ib\varphi _i^{p+1}+D^2\varphi _i+\varphi _i\eta _i\sigma $$
(1)
where $`\eta `$ is a Gaussian delta correlated noise, $`a`$, $`b`$, $`D`$, and $`\sigma `$ are constants, and $`\varphi `$ is the local order-parameter field, describing the physical density under study. In order to study its mean-field solution, we consider the case of global coupling, i.e. define the system on a fully connected (or complete) graph
$$^2\varphi _i=\frac{1}{N1}\underset{ji}{}(\varphi _j\varphi _i)M_i\varphi _i.$$
(2)
The associated Fokker-Planck equation vK for Eq. (1) in the Ito sense can then be solved in the stationary case Birner ; Walter ; Marsili , assuming $`M_i=M`$,
$$R(\varphi ;M)=\varphi ^{22(Da)/\sigma ^2}\mathrm{exp}\left[\frac{2b\varphi ^p}{p\sigma ^2}\frac{2MD}{\varphi \sigma ^2}\right].$$
(3)
This is a power-law distribution function with two cutoffs: an upper one coming from the non-linear saturation term and a lower one generated by $`M`$ which acts as a constant external field.
In order to proceed further $`M`$ is taken equal to its average over realizations, $`m`$, which is determined self-consistently by imposing vdB ; Walter ; Birner
$$m=\varphi =\frac{𝑑\varphi \varphi R(\varphi ;m)}{𝑑\varphi R(\varphi ;m)}.$$
(4)
We will denote in the following this solution, valid for $`N=\mathrm{}`$, as the Non Fluctuating (NF) solution. Equation (4) has always the trivial solution $`m=0`$, stable for $`a<a_c`$ with $`a_c=0`$. At the critical point, $`a_c`$, a stable solution with $`m>0`$ appears. Introducing the distance from the critical point $`ϵ=a/(\sigma ^2/2)`$ and the reduced variable $`s=2D/\sigma ^2`$, we can rewrite the probability distribution as
$$R(\varphi |m)=\varphi ^\alpha \mathrm{exp}\left[\frac{bs\varphi ^p}{pD}\frac{ms}{\varphi }\right]\alpha =2+sϵ.$$
(5)
Introducing the notation
$$I_{s,k}=_0^{\mathrm{}}𝑑\varphi \varphi ^kR_s(\varphi |m)$$
(6)
the first moment $`m=I_{s,1}/I_{s,0}`$ is easily computed, yielding the relation between $`m`$ and $`ϵ`$ in the active region
$$m^{\mathrm{min}[sϵ,p]}ϵ$$
(7)
so that $`mϵ^{\beta _1}`$ with $`\beta _1=\mathrm{max}[1/s,1/p]`$. In this way, two different regimes were found by Birner et al.: one is universal (weak-noise) and the other one is not (strong-noise) with the order-parameter exponent changing continuously with $`s`$. These two regimes are the analogue of the well-known weak-coupling (Gaussian) and strong-coupling (non Gaussian) regimes of KPZ dynamics HZ .
Let us now go beyond the results in Birner by computing higher moments $`m_k=\varphi ^k=I_{s,k}/I_{s,0}`$. The normalization $`I_{s,0}`$ scales as $`m^{1s+ϵ}`$, and to leading order
$$I_{s,k}=I_{sk+1,1}B_km^{s+k1+ϵ}+C_k+D_km^{ps+k1+ϵ}.$$
(8)
While the third term is always sub-leading, it depends on $`k`$ which of the other two dominates. For $`s>k1+ϵ`$ the first one is dominant, so that
$$m_k=\varphi ^k\frac{m^{k1s+ϵ}}{m^{1s+ϵ}}m^k$$
(9)
and for $`s<k1+ϵ`$ the second is the leading one, so that
$$m_k=\varphi ^k\frac{1}{m^{1s+ϵ}}m^{1+sϵ}$$
(10)
Hence, the regime with non universal first moment is quite rich. For $`1/p<1/s<1`$ the first $`k`$ moments exhibit standard scaling, as long as $`k<s+1ϵ`$, while all others scale with the same exponent $`1/s+1`$. We call this regime intermediate-noise; it does not have a clear analogue in free KPZ-like interfaces.
When $`1/s>1`$ full multi-scaling occurs: all moments, except the first, scale with the same exponent, $`1/s+1`$. This is the strong-noise regime; numerical evidence of this multi-scaling is provided in Fig. 1, and is similar to the known phenomenology of single-site (zero-dimensional) MN equations Graham ; MN .
¿From naive power counting (usually expected to reproduce mean-field exponents) performed on the MN equation, the naive time scale $`T`$ should scale as $`T^1aϵ`$ and, therefore, we expect the critical time-decay exponent of the order parameter, $`\theta _1`$ (or simply $`\theta `$) defined by $`(\varphi (t,ϵ=0)t^{\theta _1}`$) to coincide with $`\beta _1`$ in mean-field. The same property remains valid for higher moments, so the multi-scaling property of the static exponent is translated into multi-scaling of the decay exponents footnote .
Let us now construct a Ginzburg criterion within the fully connected lattice, to see under which circumstances the previous approximation ceases to be sound. To do so we just need to compute the ratio $`\frac{rm_2}{m_1^2}`$ Birgenau ; LB . Whenever $`r`$ diverges fluctuations are expected to play a significant role, breaking down the self-consistent solution which ignores them, as soon as they are taken into consideration. From the previously reported scaling for the strong-noise solution one has $`rϵ^{11/s}`$, which diverges at the critical point if $`1/s>1`$, i. e. in any case cite GGC .
Therefore, the previously reported strong-noise mean-field solution is fully valid only in the strict $`N=\mathrm{}`$ limit, when fluctuations can be safely discarded owing to the law of large numbers. In general, it is expected that such a strong-noise solution will break down as soon as fluctuations are somehow taken into account. If, for example, the MN equation is defined on the top of a d-dimensional lattice (every site sees a finite number of others) the NF strong-noise solution is not expected a priori to be valid. We will discuss this aspect more in detail in the forthcoming sections. On the other hand, for the weak- and the intermediate-noise regimes, we have Gaussian scaling of the lowest moments, and therefore $`r`$ converges to a constant at the critical point.
## III Fluctuating Solution
In the previously reported approach the crucial step is the replacement of $`M`$ (the mean value seen by any arbitrary single site) by a fixed non-fluctuating $`m`$. As said before, owing to the law of large numbers this is exact if $`N=\mathrm{}`$ whatever the probability distribution of the neighboring sites, but one can wonder whether this replacement is acceptable if a finite system size $`N`$ is considered. In other words: what is the critical behavior for large but finite values of $`N`$? This is a perfectly sensible question, since for MN a sharp phase transition is well defined for any value of $`N`$, even for $`N=1`$ Graham .
In this section we relax the condition of a fixed value for $`M`$, by allowing the spatial average value of the $`\varphi `$ field in the stationary state to fluctuate in time, sampling some probability distribution $`Q(M)`$ to be determined self-consistently.
Assuming that $`M`$ changes in a time scale much larger than the characteristic time scale for $`\varphi `$, we can still solve the Fokker-Planck for a fixed value of $`M`$, and therefore Equation (3) still holds, but now it has to be interpreted as a conditional probability $`R(\varphi M)`$. The full distribution $`P(\varphi )`$ is given by the convolution:
$$P(\varphi )=_0^{\mathrm{}}𝑑MR(\varphi |M)Q(M).$$
(11)
In this case, the self-consistent Equation (4) has to be replaced by a self-consistent functional equation for $`Q(M)`$.
$$𝑑MMQ(M)=\frac{𝑑\varphi \varphi 𝑑MQ(M)R(\varphi M)}{𝑑\varphi 𝑑MQ(M)R(\varphi M)}.$$
(12)
By solving this functional equation one could obtain the full solution as in the NF case, within the slow changing $`M`$ approximation.
Instead of solving numerically such a self-consistent equation we now try to write down an evolution equation for $`M(t)`$, from which $`Q(M)`$ follows. For that, we sum Eq. (1) over $`i`$ and divide it by $`N`$, giving
$$\dot{M}=aMb\frac{1}{N}\underset{i}{}\varphi _i^{p+1}+\sigma \frac{1}{N}\underset{i}{}\varphi _i\eta _i.$$
(13)
This turns out to be an effective MN equation for the field $`M`$ in zero dimensions. Let us justify this statement. The second term on the r. h. s. is perfectly analogous to the nonlinear term in Eq. (1): it just introduces an upper cutoff $`M_u`$ in the distribution of the field $`M`$. The last term on the r. h. s. is less trivial: it can be rewritten as $`\sigma \frac{1}{N}_i\frac{\varphi _i}{M}\eta _iM`$, hence in the form of a multiplicative noise $`\eta ^{}M`$. The average value of $`\eta ^{}`$ is clearly zero, since $`\eta ^{}`$ has random sign. Its second moment, that we denote as $`\sigma ^2(N)`$, is
$$\sigma ^2(N)=\frac{\sigma ^2}{N^2}(\underset{i}{}\frac{\varphi _i}{M}\eta _i)^2.$$
(14)
The variables $`\eta _i`$ and $`\varphi _i`$ are uncorrelated. The $`\eta _i`$ are normally distributed. In practice, to compute $`\sigma ^2(N)`$ we must evaluate the sum of $`N`$ variables $`\varphi _i/M`$, distributed according to $`R(\varphi /M|M)`$ with random signs. The distribution $`R(\varphi /M|M)`$ is a power-law with exponent $`\alpha =2+sϵ`$, lower cutoff in 1 and upper cutoff in $`1/M`$.
We now discuss the properties of this solution depending on whether we are working in the strong-noise regime or not.
### III.1 Weak- and Intermediate-Noise
Let us consider first what happens when $`1/s<1`$, i.e. in the weak and the intermediate noise regimes. In these cases, the exponent $`\alpha `$ of the distribution $`R`$ is larger than $`3`$, for $`a`$ sufficiently small. Adding $`N`$ such variables with random signs is equivalent to adding Gaussian variables Levy . Then
$$\left(\underset{i}{}\frac{\varphi _i}{M}\eta _i\right)^2N,$$
(15)
so that
$$\sigma ^2(N)\sigma ^2/N.$$
(16)
We can conclude that the dynamics of $`M`$ is governed by an effective equation with MN in zero dimensions with a noise $`\eta ^{}`$ with renormalized variance $`\sigma ^2/N`$. It is clear that $`\eta ^{}`$ is correlated in time, but its finite correlation-time can be eliminated by suitably rescaling the time variable. The distribution for $`M`$ is then given by the solution of the zero-dimensional MN Graham
$$Q(M)=M^{2[1a/\sigma ^2(N)]}\mathrm{exp}(M/M_u)^2.$$
(17)
As expected, this distribution function is a solution of the previously written functional self-consistent Equation (12).
For finite $`N`$ the system undergoes an absorbing phase transition for a finite value of the control parameter
$$a_c(N)=\frac{\sigma (N)^2}{2}\frac{\sigma ^2}{2N}.$$
(18)
The exponents for the transition with finite $`N`$ are given by the values for zero-dimensional MN transition Graham ; MN : $`\beta _1=1`$, and $`\theta =1/2`$.
Notice that this transition is qualitatively different from that occurring in the thermodynamic limit. Here the transition takes place because the distribution of $`M`$ becomes non-normalizable due to the divergence for $`M0`$. In the thermodynamic limit, instead, the distribution $`Q(M)`$ remains normalizable (Gaussian) at the transition and criticality comes from the peak position $`M_u`$ moving toward zero (see the appendix, where a coherent general scaling picture is presented).
These two regimes are therefore distinguished by the presence of a (Gaussian) peak for finite $`M`$ (the NF limit) or a broad distribution with a power-law divergence for $`M0`$. The crossover occurs where the power-law exponent in Eq. (17) changes sign, i.e. for
$$\stackrel{~}{a}(N)=\sigma ^2(N)=2a_c(N).$$
(19)
The behavior for fixed $`N`$ can then be summarized as follows. For $`a\stackrel{~}{a}(N)`$ the system exhibits the critical behavior of the weak noise $`N=\mathrm{}`$ solution. For $`0<aa_c(N)a_c(N)`$ it behaves as it was zero-dimensional. In the latter case, the crossover can also be observed with fixed $`a`$ by looking at the temporal evolution of the first moment $`m`$: at short times the NF solution is followed, crossing-over later to the asymptotic zero-dimensional scaling.
We have checked the correctness of this scenario by means of numerical simulations of Eq. (1) with $`p=2`$. From the numerical point of view, the first problem is the determination of the critical point $`a_c(N)`$. The criterion we have chosen is based on the way the first moment $`m(t)`$ decays in time. $`a_c(N)`$ is the value separating a concave behavior (for $`a>a_c(N)`$) from a convex one (for $`a<a_c(N)`$). In this way, we obtain the values plotted in Fig. 2. The behavior found in the weak-noise regime is in perfect agreement with the predicted $`1/N`$ behavior.
With $`p=2`$, the temporal behavior in the weak-noise case is the same ($`\theta =1/2`$) both in the thermodynamical limit and for the effective zero-dimensional behavior valid at finite $`N`$, not allowing to distinguish between the two regimes. In Fig. 3 we analyze the behavior of the stationary value of the first moment $`m`$. In the main part of the figure we observe that, for $`N=5000`$, the system follows very accurately (for the values of $`a`$ considered) the decay with $`\beta _1=1/2`$, expected in the thermodynamic limit. For $`N=10`$, instead, a singularity is observed for finite $`a`$. If we plot $`m`$ as a function of $`aa_c(N)`$ (Fig. 3, inset) both behaviors can be observed: for $`aa_c(N)a_c(N)0.014`$ the scaling is of the NF type ($`\beta _1=1/2`$). For $`aa_c(N)a_c(N)`$ we observe the zero-dimensional exponent $`\beta _1=1`$.
A direct validation of Eq. (17) is provided by Fig. 4, where the distribution of the self-consistent field $`M`$ exhibits a peak in $`M_u(a)`$ for $`aa_c(N)a_c(N)`$ while a power-law divergence at zero develops for $`aa_c(N)a_c(N)`$.
### III.2 Strong-Noise
Let us consider now the strong-noise case $`1/s>1`$. The general picture is similar to the one previously described, with the main difference that the exponent $`\alpha `$ in the distribution $`R(\varphi /M|M)`$ is now smaller than $`3`$. From the theory of Levy-stable distributions Levy we know that the sum of $`N`$ variables distributed as a power-law with exponent $`1+\mu `$ with $`1<\mu <2`$ and an upper cutoff $`1/M`$, scales as $`N^{1/\mu }`$ for $`NM^\mu `$ while it behaves in a Gaussian way, $`N^{1/2}`$, for $`NM^\mu `$. This implies that, for fixed $`N`$, $`\eta ^{}`$ is a power-law distributed noise with exponent $`1+\mu `$ for $`MN^{1/\mu }`$, and a Gaussian noise with $`\sigma ^2(N)\sigma ^2/N`$ for $`MN^{1/\mu }`$. In the present case $`\mu =\alpha 1=1+s2a/\sigma ^2`$.
Eq. (13) describes now a zero-dimensional MN with a rather exotic noise, whose distribution depends on $`M`$. We have no clear theoretical understanding of such a model. In principle it could give rise to completely new and non-trivial critical features. However, as shown numerically below, it turns out to behave asymptotically as the zero-dimensional case with standard noise. The only change is in the position of the critical point $`a_c(N)`$. Accordingly, we assume, rather crudely, that the power-law noise does not change the behavior of the zero-dimensional MN, except for the form of $`\sigma ^2(N)`$
$$\sigma ^2(N)\sigma ^2N^{2(1/(\alpha 1)1)}=\sigma ^2N^{2[1/(1+s2a/\sigma ^2)1]}.$$
(20)
The distribution of $`M`$ is then given again by Eq. (17), with the additional complication that $`\sigma ^2(N)`$ depends on $`a`$. The critical point is determined by the implicit condition
$$a_c(N)=\frac{\sigma ^2(N)}{2}\sigma ^2N^{2\left[{\displaystyle \frac{1}{(1+s2a_c(N)/\sigma ^2)}}1\right]}.$$
(21)
If, as a first approximation, we neglect the dependence of $`\alpha `$ on $`a`$, i.e. we take $`\alpha =2+s`$, we obtain $`a_c(N)N^{2s/(1+s)}`$, not far from the numerical results of Fig. 2.
Again for $`a>\stackrel{~}{a}(N)=2a_c(N)`$ the exponent of the distribution $`Q(M)`$ becomes negative and there is a crossover to the NF limit. In Fig. 5 we plot $`m`$ versus $`a`$ in the strong-noise regime. Again the value $`\beta _1=1/s`$, valid in the thermodynamic limit, is observed for large $`N`$, as the critical point and the crossover are very close to $`0`$. Near the transition the exponent is the zero-dimensional one, $`\beta _1=1`$, indicating that the non trivial noise does not modify the zero-dimensional critical behavior.
As in the other case, it is interesting to look also at the temporal evolution of $`m`$. In the strong-noise regime, since the exponent $`\theta `$ is different in the NF solution and in the zero-dimensional regime, the crossover between the two regimes results in a crossover in the decay of $`m(t)`$. This is evident from Fig. 6. The effective exponent switches from a value close to the prediction $`\theta =1/s=2`$ for short times, to the zero-dimensional value $`\theta =1/2`$, for longer times.
### III.3 Discussion
Let us underline the non-commutativity of the limits $`aa_c`$ and $`N\mathrm{}`$. It is only when the thermodynamic limit is taken first and a homogeneous mean-field $`M`$ is considered, that the NF solution is recovered. On the other hand, if one takes first the limit $`aa_c(N)`$ for a generic finite value of $`N`$, the zero-dimensional solution always dominates the scaling nearby the critical point, no matters how large $`N`$. Therefore, the inclusion of fluctuations has a dramatic effect on the NF solution.
Such a conclusion holds for all the regimes considered: weak, intermediate, and strong-noise. This is due to the fact that the two aforementioned limits do not commute, implying the presence of a non-analyticity of the most general solution in a neighborhood of the critical point. This is analogous to the observation of Gaussian scaling in standard phase transitions whenever the system size is not infinite; it is only when the thermodynamic limit is taken that the true asymptotic scaling emerges. In the case studied here the role of the Gaussian scaling is replaced by a MN zero-dimensional non-trivial scaling. This breaking down of the thermodynamic-limit behavior for finite systems has, in principle, nothing to do with the previously constructed Ginzburg criterion which leads to a breakdown of the mean-field solution only at the strong-noise regime.
In the appendix we present a general scaling theory accounting in a compact form for all the previously discussed phenomenology.
## IV Finite connectivity: Random Regular Graphs
In order to shed some light on the question of whether the NF solution holds or not for an arbitrarily large space dimensionality, $`d`$, in which the number of sites seen by any given one is finite ($`2d`$ for a hyper-cubic lattice) we should first answer the following question: does the NF behavior emerge because the size of the system goes to infinity or because the number of nearest neighbors diverges? (Note that in the fully connected graph these two limits coincide). In order to clarify this point we have considered the MN on a connected Regular Random Graph (RRG), where each site has fixed degree $`k>2`$ and random connections (for $`k=2`$ we have a one-dimensional lattice). In this case, as the number of neighbors is finite for each site, we expect fluctuations in $`M`$ and therefore a possible breakdown of the NF solution (at least in the strong-noise limit), even if the large system size limit is taken.
However, numerical results disprove such a conjecture, as we show in what follows. We have performed simulations of a system with $`k=10`$ and growing $`N`$. It turns out that the position of the critical point $`a_c(N)`$ depends on $`N`$ and does not reach a finite value dependent only on $`k`$ (see Fig. 7). Its behavior is not very different from what occurs on the fully connected system (see Eq. 21).
Moreover, one can monitor the temporal behavior of $`m(t)`$ for $`aa_c(N)`$ (Fig. 8). One finds a crossover from a NF behavior at short time to a zero-dimensional behavior at longer time, exactly as for the fully connected graph.
We conclude that the thermodynamical limit is described, also for a regular random graph with finite connectivity, by the NF behavior. The zero-dimensional behavior holds only as long as criticality is studied for finite system sizes. A discussion of these facts is presented in the next section.
## V Summary and Conclusions
In this paper we have investigated the properties of the mean-field solution for systems with multiplicative noise. First, we have revisited the self-consistent solution obtained in Ref. Birner by assuming that $`M`$ is a homogeneous non-fluctuating field. We have shown that three regimes can actually be identified depending on the noise amplitude $`1/s=\sigma ^2/(2D)`$. In the weak-noise case ($`1/s<1/p`$) all moments of the order parameter distribution obey ordinary Gaussian scaling. In the intermediate regime, ($`1/p<1/s<1`$), Gaussian scaling holds only up to a certain order, with multi-scaling emerging for higher moments. In the strong-noise regime ($`1/s>1`$) all moments scale with exponent $`1/s+1`$ except for the first one that goes as $`ϵ^{1/s}`$.
Complete absence of fluctuations for the spatial average $`M`$ holds only in the thermodynamic limit, with a strictly infinite number of nearest neighbors. In this paper we have relaxed the condition that $`M`$ can assume only one (self-consistently fixed) value and let it fluctuate: this is equivalent to considering a fully connected graph with a finite number $`N`$ of nodes.
When $`N`$ is finite $`M`$ fluctuates and the critical behavior is not described anymore by the NF thermodynamic limit, in none of the three previously described regimes. While in the NF case, criticality is governed by fluctuations of the $`\varphi `$ field in a single realization around its mean value, for finite $`N`$ what actually matters are the fluctuations of the self-consistent field $`M`$. They are described by an effective equation for a zero-dimensional system with multiplicative noise, which exhibits a critical behavior distinct from the one displayed in the thermodynamic limit.
The crossover between the two types of behavior is well described by the zero-dimensional description: both the crossover point and the critical point location depend on $`N`$ so that for any finite value of $`N`$ there is a finite interval of values of the control parameter where the effective zero-dimensional behavior is observed.
Finally, by means of numerical simulations on connected Random Regular Graphs, we have also shown that all these results still hold (at least qualitatively) for random regular graphs with finite connectivity. This is somehow counterintuitive, and seems to contradict the previously introduced Ginzburg criterion for the strong-noise regime. Indeed, given that the fluctuations of $`M`$ are predicted by the Ginzburg criterion to be relevant in the strong-noise regime of the NF solution when the number of nearest neighbors is not infinity, it is hard to understand, why the solution in the random regular graph, in which $`M_i`$ is the fluctuating average taken over the $`k`$ neighbors of any given site $`i`$, behaves so similarly (for any value of $`N`$) to the solution in the fully connected network in which $`M=M_i`$ is fluctuation-less.
We do not have a clear explanation of this fact, but we believe that this is so because of the small-world SW nature of the RRG topology. The small-world property implies that one can reach any arbitrary node starting from a generic site with a small number of steps following network links. In this way, every site is nearby any other one, making it difficult to create “local patches” with an over-density or sub-density, which would give rise to inhomogeneities and a broad field distribution. At this point, it would be interesting to study the MN equation on a Cayley tree to see if, by introducing well defined spatial neighborhoods, the previous results and interpretation are sustained. We are presently analyzing such a problem, but prefer to leave the delicate issues involved in such a study for a future publication.
Many interesting questions remain to be answered. The main one is whether the obtained NF solution in the strong-noise regime is valid in arbitrarily large but finite physical dimensions, or whether it emerges only in $`d=\mathrm{}`$ (which is obviously related to the existence of a finite upper critical dimension for MN and KPZ problems). The Ginzburg criterion for the strong-noise solution seems to point out to the second possibility, while the fact that in the finite-connectivity random regular graphs the NF solution emerges in the thermodynamic limit, could be interpreted as supporting the first one. Further analysis along these lines is left for future work.
It remains also to be understood which is the fate of the intermediate-noise regime once the MN equation is embedded into a $`d`$-dimensional lattice. Does it simply disappear? Does it survive, introducing some type of anomalous effect for high-order moments? Does it have any analogous in KPZ-like systems?
It is our hope that this work will stimulate future research in this exciting and ever surprising field of systems with multiplicative noise.
## Acknowledgments
We acknowledge useful discussions with L. Pietronero, and P. L. Garrido. We specially thank A. Gabrielli who participated in the early stages of this project, and G. Parisi for useful comments on the first version of the manuscript. M. A. M. acknowledges financial support from the Spanish MCyT (FEDER) under project BFM2001-2841, and from the Acción Integrada Hispano-Italiana HI2003-0344.
## Appendix: General scaling of the distributions
It is possible to summarize all the scaling regimes described in the paper in a general scaling form of the distributions. This describes the crossover between the zero-dimensional case and the thermodynamic limit as $`N`$ grows. The equations obeyed by the distribution $`R(\varphi |M;a,N)`$ of the $`\varphi `$ field conditional to a value $`M`$ and by the full distribution $`P(\varphi ;a,N)`$ are
$$\{\begin{array}{ccc}R(\varphi |M;a)\hfill & & \varphi ^\alpha \mathrm{exp}\left(\frac{bs\varphi ^p}{pD}\frac{Ms}{\varphi }\right)\hfill \\ P(\varphi ;a,N)\hfill & =& _0^{\mathrm{}}𝑑MR(\varphi |M;a)Q(M;a,N).\hfill \end{array}$$
(22)
The solution of these equations found assuming no fluctuations for $`M`$ is
$$\{\begin{array}{ccc}R(\varphi ;M)\hfill & & \varphi ^\alpha \mathrm{exp}\left(\frac{bs\varphi ^p}{pD}\frac{Ms}{\varphi }\right)\hfill \\ Q(M;a,N)\hfill & & \delta [Mm(a)]\hfill \\ P(\varphi ;a,N)\hfill & =& R(\varphi |M)|_{M=m(a)}\hfill \end{array}$$
(23)
A more general explicit ansatz for the solution to this set of equations, valid for generic $`N`$, is
$$\{\begin{array}{ccc}R(\varphi |M)\hfill & & M^{\alpha 1}\varphi ^\alpha \mathrm{exp}\left(\frac{bs\varphi ^p}{pD}\frac{Ms}{\varphi }\right)\hfill \\ Q(M;a,N)\hfill & & M^\alpha ^{}e^{\left(\frac{M}{M_u}\right)^2}\hfill \\ P(\varphi ;a,N)\hfill & & \{\begin{array}{cc}\frac{\varphi ^\alpha ^{}}{\alpha \alpha ^{}}\hfill & \varphi <M_u(a)\hfill \\ \frac{\varphi ^\alpha e^{\frac{bs\varphi ^p}{pD}}}{\alpha \alpha ^{}}M_u^{\alpha \alpha ^{}}\hfill & \varphi >M_u(a)\hfill \end{array}\hfill \end{array}$$
(24)
where $`\alpha ^{}=2[1a/\sigma ^2(N)]`$. The form of $`Q(M)`$ encodes the different critical behaviors. In the NF limit ($`N`$ going to $`\mathrm{}`$ first) $`\alpha ^{}`$ is negative. In this way $`Q(M)`$ tends to $`\delta (MM_u)`$; moments go to zero because $`M_u`$ goes to zero as $`a^\beta `$. In the opposite nontrivial limit ($`ϵ`$ going to zero with $`N`$ fixed) $`\alpha ^{}`$ is positive and grows up to $`1`$ for $`aa_c(N)`$. In this case, the moments go to zero because the distribution $`Q(M)`$ develops a peak in zero. The critical behavior of the moments in this limit is governed by the way $`\alpha ^{}`$ goes to $`1`$. The origin of the two limits ($`ϵ0`$ and $`N\mathrm{}`$) non-commutativity has its roots in the non-analytical form of $`Q(M`$ at $`M=0`$. |
warning/0506/math0506472.html | ar5iv | text | # Contact Dehn surgery, symplectic fillings, and Property P for knots
## 1 Property P for knots
According to a fundamental theorem of Lickorish and Wallace from the 1960s, every closed, connected, orientable $`3`$–manifold can be obtained by performing Dehn surgery on a link in the $`3`$–sphere. Previous to the recent work of Perelman, which is expected to close the coffin on the Poincaré conjecture, it was a natural question for geometric topologists whether one might be able to produce a counterexample to that conjecture by a single Dehn surgery. This led to the definition of the following property, whose name is generally regarded as a little unfortunate.
###### Definition.
A knot $`K`$ in $`S^3`$ has Property P if every nontrivial surgery along $`K`$ yields a non-simply-connected $`3`$–manifold.
Our knots are always understood to be smooth, or at least tame, i.e. equivalent to a smooth one.
Let me briefly recall the notion of Dehn surgery along a knot $`K`$ in the $`3`$–sphere $`S^3`$. Write $`\nu KS^1\times D^2`$ for a (closed) tubular neighbourhood of $`K`$. On the boundary $`(\nu K)T^2`$ of this tubular neighbourhood there are two distinguished curves (which we implicitly identify with the classes they represent in the homology group $`H_1(T^2)`$):
* The meridian $`\mu `$, defined as a simple closed curve that generates the kernel of the homomorphism on $`H_1`$ induced by the inclusion $`T^2\nu K`$.
* The preferred longitude $`\lambda `$, defined as a simple closed curve that generates the kernel of the homomorphism on $`H_1`$ induced by the inclusion $`T^2C:=\overline{S^3\nu K}`$.
This preferred longitude can also be characterised by the property that it has linking number zero with $`K`$. The knot $`K`$ bounds an embedded surface in $`S^3`$ (called a Seifert surface for $`K`$), and $`\lambda `$ can be obtained by pushing $`K`$ along that surface. For that reason, the trivialisation of the normal bundle of $`K`$ defined by $`\lambda `$ is called the surface framing of $`K`$.
Given an orientation of $`S^3`$, orientations of $`\mu `$ and $`\lambda `$ are chosen such that $`(\mu ,\lambda )`$ is a positive basis for $`H_1(T^2)`$, with $`T^2`$ oriented as the boundary of $`\nu K`$. In the contact geometric setting below, the orientation of $`S^3`$ will be the one induced from the contact structure.
Let $`p,q`$ be coprime integers. The manifold $`K_{p/q}`$ obtained from $`S^3`$ by Dehn surgery along $`K`$ with surgery coefficient $`p/q\{\mathrm{}\}`$ is defined as
$$K_{p/q}:=\overline{S^3\nu K}_gS^1\times D^2,$$
where the gluing map $`g`$ sends the meridian $`\times D^2`$ to $`p\mu +q\lambda `$. The resulting manifold is completely determined by the knot and the surgery coefficient.
A simple Mayer-Vietoris argument shows that $`H_1(K_{p/q})_{|p|}`$. Therefore, saying that a knot $`K`$ has Property P is equivalent to
$$\pi _1(K_{1/q})=1\text{only for}q=0.$$
(Observe that $`p/q=\mathrm{}`$ corresponds to a trivial surgery.)
###### Example.
The unknot does not have Property P. Indeed, every $`(1/q)`$–surgery on the unknot yields $`S^3`$, which is seen as follows. If $`K`$ is the unknot, then the closure $`C`$ of $`S^3\nu K`$ is also a solid torus. Write $`\mu _C`$ and $`\lambda _C`$ for meridian and preferred longitude on $`C`$. We may assume $`\mu =\lambda _C`$ and $`\lambda =\mu _C`$. When performing $`(1/q)`$–surgery on $`K`$, a solid torus is glued to $`C`$ by sending its meridian $`\mu _0`$ to $`\mu +q\lambda =\lambda _C+q\mu _C`$. Now, there clearly is a diffeomorphism of $`C`$ that sends $`\mu _C`$ to itself and $`\lambda _C`$ to $`\lambda _C+q\mu _C`$. It follows that the described surgery is equivalent to the one where we send $`\mu _0`$ to $`\lambda _C=\mu `$, which is a trivial $`\mathrm{}`$–surgery.
In the early 1970s, Bing and Martin, as well as González-Acuña, conjectured that every nontrivial knot has Property P. By work of Kronheimer and Mrowka , this is now a theorem.
###### Theorem 1 (Kronheimer-Mrowka).
Every nontrivial knot in $`S^3`$ has Property P.
Before describing the role that contact geometry has played in the proof of this theorem, I want to indicate the importance of this theorem beyond the negative statement that counterexamples to the Poincaré conjecture cannot result from a single surgery.
###### Proposition 2.
If two knots $`K,K^{}`$ in $`S^3`$ have homeomorphic complements and one of the knots has property P, then the knots are equivalent, i.e. there is a homeomorphism of $`S^3`$ mapping $`K`$ to $`K^{}`$.
###### Proof.
According to a result of Edwards , two compact $`3`$–manifolds with boundary are homeomorphic if and only if their interiors are homeomorphic. Thus, if $`S^3K`$ is homeomorphic to $`S^3K^{}`$, then there is a homeomorphism $`\phi :CC^{}`$, where $`C:=\overline{S^3\nu K}`$ and $`C^{}:=\overline{S^3\nu K^{}}`$.
Suppose $`K`$ has Property P. This implies that there is a unique way of attaching a solid torus $`S^1\times D^2`$ to $`C`$ such that the resulting manifold is the $`3`$–sphere. Hence $`\phi `$ extends to a homeomorphism $`S^3S^3`$, i.e. the knots $`K`$ and $`K^{}`$ are equivalent. ∎
Observe that in this proof we only used the weaker property that nontrivial surgery along $`K`$ does not yield the standard $`3`$–sphere. This had been proved earlier (for $`K`$ different from the unknot) by Gordon and Luecke . Since the unknot is characterised by its complement being a solid torus, the result of Kronheimer and Mrowka (or the weaker one by Gordon and Luecke) yields the following corollary.
###### Corollary 3.
If two knots in $`S^3`$ have homeomorphic complements, then the knots are equivalent. ∎
Of course, together with a positive answer to the Poincaré conjecture, the result of Gordon-Luecke implies that of Kronheimer-Mrowka.
## 2 Contact Dehn surgery
This section gives a brief report on joint work with Fan Ding . Recall that a (coorientable) contact structure $`\xi `$ on a differential $`3`$–manifold is a tangent $`2`$–plane field defined as the kernel of a global differential $`1`$–form $`\alpha `$ that satisfies the nonintegrability condition $`\alpha d\alpha 0`$ (meaning that $`\alpha d\alpha `$ vanishes nowhere). An example is the standard contact structure
$$\xi _{st}=\mathrm{ker}(xdyydx+zdttdz)$$
on $`S^3^4`$. This can also be characterised as the complex line in the tangent bundle $`TS^3`$ of $`S^3`$ with respect to complex multiplication induced from the inclusion $`TS^3T^2|_{S^3}`$.
I shall have to use a few notions from contact geometry without time for much explanation (tight and overtwisted contact structures, convex surfaces in contact $`3`$–manifolds). For more details see the introductory lectures by Etnyre or the Handbook chapter by the present author .
A (smooth) knot $`K`$ in a contact $`3`$–manifold $`(M,\xi )`$ is called Legendrian if it is everywhere tangent to $`\xi `$. The normal bundle of such a knot has a canonical trivialisation, determined by a vector field along $`K`$ that is everywhere transverse to $`\xi `$. This will be referred to as the contact framing. We now consider Dehn surgery along $`K`$ with coefficient $`p/q`$ as before, but we define the surgery coefficient with respect to the contact framing.
It turns out that for $`p0`$ one can always extend the contact structure $`\xi |_{M\nu K}`$ to one on the surgered manifold in such a way that the extended contact structure is tight on the glued-in solid torus $`S^1\times D^2`$. Moreover, subject to this tightness condition there are but finitely many choices for such an extension, and for $`p/q=1/k`$ with $`k`$ the extension is in fact unique. These observations hinge on the fact that $`(\nu K)`$ is a convex surface, i.e. a surface admitting a transverse flow preserving the contact structure. On solid tori with convex boundary condition, tight contact structures have been classified by Giroux and Honda. Furthermore, one knows how to glue contact manifolds along convex surfaces, since the germ of a contact structure along a convex surface is determined by some simple data on that surface.
We can therefore speak sensibly of contact $`(1/k)`$–surgery. The following theorem is proved in .
###### Theorem 4.
Let $`(M,\xi )`$ be a closed, connected contact $`3`$–manifold. Then $`(M,\xi )`$ can be obtained from $`(S^3,\xi _{st})`$ by contact $`(\pm 1)`$–surgery along a Legendrian link. ∎
###### Remarks.
(1) There is a related theorem, due to Lutz-Martinet in the early 1970s, cf. , saying that every (closed, orientable) $`3`$–manifold admits a contact structure in each homotopy class of tangent $`2`$–plane fields. The original proof is based on surgery along a link in $`S^3`$ transverse to $`\xi _{st}`$. For an alternative proof using Legendrian surgery see .
(2) From the topological point of view, surgeries with integer surgery coefficient are best, since they correspond to attaching $`2`$–handles to the boundary of a $`4`$–manifold. Thus, contact $`(\pm 1)`$–surgeries are best from both the topological and contact geometric viewpoint.
(3) If $`(M^{},\xi ^{})`$ is obtained from $`(M,\xi )`$ by a contact $`(1/k)`$–surgery, one can recover $`(M,\xi )`$ by a suitable contact $`(1/k)`$–surgery on $`(M^{},\xi ^{})`$, see .
(4) Contact $`(1)`$–surgery is symplectic handlebody surgery in the sense of Eliashberg and Weinstein, cf. , and preserves the property of being strongly symplectically fillable (see below).
## 3 Symplectic fillings
Contact geometry enters the proof of Theorem 1 via the notion of symplectic fillings. Observe that a contact $`3`$–manifold $`(M,\xi )`$ is naturally oriented — the sign of the volume form $`\alpha d\alpha `$ does not depend on the choice of $`1`$–form $`\alpha `$ defining a given $`\xi `$; similarly, a symplectic $`4`$–manifold $`(W,\omega )`$, i.e. with $`\omega `$ a closed $`2`$–form satisfying $`\omega ^20`$, is naturally oriented by the volume form $`\omega ^2`$.
###### Definition.
(a) A compact symplectic $`4`$–manifold $`(W,\omega )`$ is called a weak (symplectic) filling of the contact manifold $`(M,\xi )`$ if $`W=M`$ as oriented manifolds (outward normal followed by orientation of $`M`$ gives orientation of $`W`$) and $`\omega |_\xi 0`$.
(b) A compact symplectic $`4`$–manifold $`(W,\omega )`$ is called a strong (symplectic) filling of the contact manifold $`(M,\xi )`$ if $`W=M`$ and there is a Liouville vector field $`X`$ defined near $`W`$, pointing outwards along $`W`$, and satisfying $`\xi =\mathrm{ker}(i_X\omega |_{TM})`$. Here Liouville vector field means that the Lie derivative $`_X\omega `$, which is the same as $`d(i_X\omega )`$ because of $`d\omega =0`$ and Cartan’s formula, is required to equal $`\omega `$.
For instance, $`(S^3,\xi _{st})`$ is strongly filled by the standard symplectic $`4`$–disc $`D^4`$ with $`\omega _{st}=dxdy+dzdt`$. The Liouville vector field here is the radial vector field $`X=r_r/2`$.
It is clear that every strong filling is also a weak filling. The converse is false: There are contact structures that are weakly but not strongly fillable; such examples are due to Eliashberg and Ding-Geiges.
The contact geometric result that allowed Kronheimer and Mrowka to conclude their proof of Property P was established by Eliashberg .
###### Theorem 5 (Eliashberg).
Any weak symplectic filling of a contact $`3`$–manifold embeds symplectically into a closed symplectic $`4`$–manifold.
An alternative proof was given by Etnyre . Both proofs rely on open book decompositions adapted (in the sense of Giroux) to contact structures. Theorem 5 being a cobordism theoretic result, it is arguably more natural to give a surgical proof. Özbağcı and Stipsicz were the first to observe that such a proof, based on Theorem 4, can indeed be devised. In the remainder of this section, I shall sketch this surgical argument.
Theorem 5 is proved by showing that any contact $`3`$–manifold has what is called a concave filling that can be glued to the given (convex) filling. (For instance, a strong concave filling corresponds to a Liouville vector field pointing inwards along the boundary.) Such a ‘cap’, attached to the (convex) symplectic filling of the contact manifold, gives the desired closed symplectic manifold.
(i) Strong fillings can be capped off: Let $`(W,\omega )`$ be a strong filling of $`(M,\xi )`$. By Theorem 4, there is a Legendrian link $`𝕃=𝕃^{}𝕃^+`$ in $`(S^3,\xi _{st})`$ such that contact $`(1)`$–surgery along the components of $`𝕃^{}`$ and contact $`(+1)`$–surgery along those of $`𝕃^+`$ produces $`(M,\xi )`$. By Remarks (3) and (4) we can attach symplectic $`2`$–handles to the boundary $`(M,\xi )`$ of $`(W,\omega )`$ corresponding to contact $`(1)`$–surgeries that undo the contact $`(+1)`$–surgeries along $`𝕃^+`$. The result will be a symplectic manifold $`(W^{},\omega ^{})`$ strongly filling a contact manifold $`(M^{},\xi ^{})`$, and the latter can be obtained from $`(S^3,\xi _{st})=(D^4,\omega _{st})`$ by performing contact $`(1)`$–surgeries (along $`𝕃^{}`$) only.
A handlebody obtained from $`(D^4,\omega _{st})`$ by attaching symplectic handles in this way is in fact a Stein filling of its boundary contact manifold, and for those a symplectic cap had been found earlier by Akbulut-Özbağcı and Lisca-Matić. The cap that fits on the Stein filling also fits on the strong filling $`(W^{},\omega ^{})`$, since strongly convex and strongly concave fillings of a given contact manifold can always be glued together, using the Liouville flow to define collar neighbourhoods of the boundary.
(ii) Reduce the problem to the consideration of homology spheres only: Let $`(W,\omega )`$ be a weak filling of $`(M,\xi )`$. We want to attach a (weak) symplectic cobordism from $`(M,\xi )`$ to some integral homology sphere $`\mathrm{\Sigma }^3`$ with contact structure $`\xi ^{}`$, so as to get a weak filling of $`(\mathrm{\Sigma }^3,\xi ^{})`$ containing $`(M,\xi )`$ as a separating hypersurface.
We start from a contact surgery presentation of $`(M,\xi )`$ as in (i). For each component $`L_i`$ of $`𝕃`$ we choose a Legendrian knot $`K_i`$ in $`(S^3,\xi _{st})`$ only linked with that component, with linking number $`1`$. These $`K_i`$ can be chosen in such a way that surgery with coefficient $`1`$ relative to the contact framing is the same as surgery with coefficient $`0`$ relative to the surface framing. (In case you know the term: The Thurston-Bennequin invariant of $`K_i`$ can be chosen to be equal to $`1`$). Performing these surgeries has the effect of killing the first integral homology.
Since $`\omega `$ is exact in the neighbourhood $`S^1\times D^2\times (\epsilon ,0]`$ of a Legendrian knot in the boundary $`(M,\xi )`$ of $`(W,\omega )`$, these surgeries can be performed by attaching symplectic $`2`$–handles as in the case of a strong filling. The collection of these handles gives the desired (weak) symplectic cobordism.
(iii) Pass from a weak filling of a homology sphere to a strong filling: We begin with a weak filling $`(W,\omega )`$ of an integral homology sphere $`(\mathrm{\Sigma }^3,\xi )`$, for instance the one obtained in (ii); beware that we retain the original notation for the filling. We want to modify $`\omega `$ in a collar neighbourhood $`\mathrm{\Sigma }^3\times [0,1]`$ of the boundary $`\mathrm{\Sigma }^3\mathrm{\Sigma }^3\times \{1\}`$ such that the resulting symplectic manifold is a strong filling of the new induced contact structure on the boundary. By (i) this can then be capped off.
Since $`H^2(\mathrm{\Sigma }^3)=0`$, we can write $`\omega =d\eta `$ with some $`1`$–form $`\eta `$ in a collar neighbourhood as described. (We see that it would be enough to have $`\mathrm{\Sigma }^3`$ a rational homology sphere.) Choose a $`1`$–form $`\alpha `$ on $`\mathrm{\Sigma }^3`$ with $`\xi =\mathrm{ker}\alpha `$ and $`\alpha \omega |_{T\mathrm{\Sigma }^3}>0`$, which is possible for a weak filling. Then set
$$\stackrel{~}{\omega }=d(f\eta )+d(g\alpha )$$
on $`\mathrm{\Sigma }^3\times [0,1]`$, where the smooth functions $`f(t)`$ and $`g(t)`$, $`t[0,1]`$, are chosen as follows: Fix a small $`\epsilon >0`$. Choose $`f:[0,1][0,1]`$ identically $`1`$ on $`[0,\epsilon ]`$ and identically $`0`$ near $`t=1`$. Choose $`g:[0,1]_0^+`$ identically $`0`$ near $`t=0`$ and with $`g^{}(t)>0`$ for $`t>\epsilon /2`$.
We compute
$$\stackrel{~}{\omega }=f^{}dt\eta +f\omega +g^{}dt\alpha +gd\alpha ,$$
whence
$`\stackrel{~}{\omega }^2`$ $`=`$ $`2ff^{}dt\eta \omega +2f^{}gdt\eta d\alpha +f^2\omega ^2`$
$`+2fg^{}\omega dt\alpha +2fg\omega d\alpha +2gg^{}dt\alpha d\alpha .`$
The terms appearing with the factors $`f^2`$, $`fg^{}`$ and $`gg^{}`$ are positive volume forms. By choosing $`g`$ small on $`[0,\epsilon ]`$ and $`g^{}`$ large compared with $`|f^{}|`$ and $`g`$ on $`[\epsilon ,1]`$, one can ensure that these positive terms dominate the three terms we cannot control. Then $`\stackrel{~}{\omega }`$ is a symplectic form on the collar and, in terms of the coordinate $`s:=\mathrm{log}g(t)`$, the symplectic form looks like $`d(e^s\alpha )`$ near the boundary, with Liouville vector field $`_s`$.
## 4 Proof of Property P for nontrivial knots
Here is a very rough sketch of the proof by Kronheimer and Mrowka. It relies heavily on pretty much everything known under the sun about gauge theory.
Let $`K`$ be a nontrivial knot. It had been proved earlier by Culler-Gordon-Luecke-Shalen that $`\pi _1(K_{1/q})`$ is nontrivial for $`q\{0,\pm 1\}`$. It therefore suffices to find a nontrivial homomorphism $`\pi _1(K_1)\text{SO}(3)`$.
Arguing by contradiction, we assume that no such homomorphism exists. This implies the vanishing of the instanton Floer homology group $`HF(K_1)`$. By the Floer exact triangle one finds that the group $`HF(K_0)`$ vanishes likewise, and so does the Fukaya-Floer homology group.
For $`K`$ nontrivial, results of Gabai say that $`K_0`$ is different from $`S^1\times S^2`$ and admits a taut $`2`$–dimensional foliation. Eliashberg and Thurston, in their theory of confoliations, deduce from this the existence of a symplectic structure on $`K_0\times [1,1]`$ weakly filling contact structures on the boundary components. According to Theorem 5, by capping off these boundaries we find a closed symplectic $`4`$–manifold $`V`$ containing $`K_0`$ as a separating hypersurface (and satisfying some mild cohomological conditions).
Now, on the one hand, the Donaldson invariants of $`V`$ can be expressed as a pairing on the Fukaya-Floer homology group of $`K_0`$ and therefore have to vanish.
On the other hand, results of Taubes say that the Seiberg-Witten invariants of $`V`$ are nontrivial. By a conjecture of Witten, proved in the relevant case by Feehan-Leness, the Donaldson invariants are likewise nontrivial. This contradiction proves Theorem 1. |
warning/0506/cs0506096.html | ar5iv | text | # Polynomial Synthesis of Asynchronous Automata
## Introduction
One of the major contributions in the theory of Mazurkiewicz traces characterizes regular languages by means of asynchronous automata which are devices with a distributed control structure. So far all known constructions of asynchronous automata from regular trace languages are quite involved and yield an exponential explosion of the number of states . Furthermore conversions of non-deterministic asynchronous automata into deterministic ones rely on Zielonka’s time-stamping function and suffer from the same state-explosion problem. Interestingly heuristics to build small deterministic asynchronous automata were proposed recently in .
Zielonka’s theorem and related techniques are fundamental tools in concurrency theory. For instance they are useful to compare the expressive power of classical models of concurrency such as Petri nets, asynchronous systems, and concurrent automata . These methods have been adapted already to the construction of communicating finite-state machines from regular sets of message sequence charts . More recently the construction of asynchronous cellular automata was used to implement globally-cooperative high-level message sequence charts . All these constructions yield an exponential explosion of the number of local states.
In this paper we give a *polynomial* construction of non-deterministic asynchronous automata. Our algorithm starts from the specification of a regular trace language in the form of a possibly non-deterministic automaton. The latter is unfolded inductively on the alphabet into an automaton that enjoys several structural properties (Section 2). Next the unfolding automaton is used as the common skeleton of all local processes (Subsection 3.2). Our algorithm is designed specifically to ensure that the number of local states built is polynomial in the number of global states in the specification (Subsection 3.1). We show how this approach subsumes the complexity of Zielonka’s and Pighizzini’s constructions (Subsection 1.3).
## 1 Background and main result
In this paper we fix a finite alphabet $`\mathrm{\Sigma }`$ provided with a total order $``$. An automaton over a subset $`T\mathrm{\Sigma }`$ is a structure $`𝒜=(Q,ı,T,\stackrel{}{},F)`$ where $`Q`$ is a *finite* set of states, $`ıQ`$ is an initial state, $`\stackrel{}{}Q\times T\times Q`$ is a set of transitions, and $`FQ`$ is a subset of final states. We write $`q\stackrel{a}{}q^{}`$ to denote $`(q,a,q^{})\stackrel{}{}`$. Then the automaton $`𝒜`$ is called *deterministic* if we have $`q\stackrel{a}{}q^{}q\stackrel{a}{}q^{\prime \prime }q^{}=q^{\prime \prime }`$. For any word $`u=a_1\mathrm{}a_n\mathrm{\Sigma }^{}`$, we write $`q\stackrel{u}{}q^{}`$ if there are some states $`q_0,q_1,\mathrm{},q_nQ`$ such that $`q=q_0\stackrel{a_1}{}q_1\mathrm{}q_{n1}\stackrel{a_n}{}q_n=q^{}`$. A state $`qQ`$ is *reachable* if $`ı\stackrel{u}{}q`$ for some $`u\mathrm{\Sigma }^{}`$. The language $`L(𝒜)`$ accepted by some automaton $`𝒜`$ consists of all words $`u\mathrm{\Sigma }^{}`$ such that $`ı\stackrel{u}{}q`$ for some $`qF`$. A subset of words $`L\mathrm{\Sigma }^{}`$ is *regular* if it is accepted by some automaton.
### 1.1 Mazurkiewicz traces
We fix an *independence relation* $``$ over $`\mathrm{\Sigma }`$, that is, a binary relation $`\mathrm{\Sigma }\times \mathrm{\Sigma }`$ which is irreflexive and symmetric. For any subset of actions $`T\mathrm{\Sigma }`$, the *dependence graph* of $`T`$ is the undirected graph $`(V,E)`$ whose set of vertices is $`V=T`$ and whose edges denote dependence, i.e. $`\{a,b\}Ea\overline{)}b`$.
The *trace equivalence* $``$ associated with the independence alphabet $`(\mathrm{\Sigma },)`$ is the least congruence over $`\mathrm{\Sigma }^{}`$ such that $`abba`$ for all pairs of independent actions $`ab`$. For a word $`u\mathrm{\Sigma }^{}`$, the *trace* $`[u]=\{v\mathrm{\Sigma }^{}|vu\}`$ collects all words that are equivalent to $`u`$. We extend this notation from words to sets of words in a natural way: For all $`L\mathrm{\Sigma }^{}`$, we put $`[L]=\{v\mathrm{\Sigma }^{}|uL,vu\}`$.
A *trace language* is a subset of words $`L\mathrm{\Sigma }^{}`$ that is closed for trace equivalence: $`uLvuvL`$. Equivalently we require that $`L=[L]`$. With no surprise a trace language $`L`$ is called *regular* if it is accepted by some automaton.
### 1.2 Asynchronous systems vs. asynchronous automata
Two classical automata-based models are known to correspond to regular trace languages. Let us first recall the basic notion of asynchronous systems .
###### Definition 1.1
An automaton $`𝒜=(Q,ı,\mathrm{\Sigma },\stackrel{}{},F)`$ over the alphabet $`\mathrm{\Sigma }`$ is called an *asynchronous system* over $`(\mathrm{\Sigma },)`$ if we have
* $`q_1\stackrel{a}{}q_2q_2\stackrel{b}{}q_3ab`$ implies $`q_1\stackrel{b}{}q_4q_4\stackrel{a}{}q_3`$ for some $`q_4Q`$.
The Independent Diamond property ID ensures that the language $`L(𝒜)`$ of any asynchronous system is closed for the commutation of independent adjacent actions. Thus it is a regular trace language. Conversely it is easy to observe that *any regular trace language is the language of some deterministic asynchronous system.*
We recall now a more involved model of communicating processes known as asynchronous automata . A finite family $`\delta =\left(\mathrm{\Sigma }_k\right)_{kK}`$ of subsets of $`\mathrm{\Sigma }`$ is called a *distribution of $`(\mathrm{\Sigma },)`$* if we have $`a\overline{)}bkK,\{a,b\}\mathrm{\Sigma }_k`$ for all actions $`a,b\mathrm{\Sigma }`$. Note that each subset $`\mathrm{\Sigma }_k`$ is a clique of the dependence graph $`(\mathrm{\Sigma },\overline{)})`$ and a distribution $`\delta `$ is simply a clique covering of $`(\mathrm{\Sigma },\overline{)})`$. We fix an arbitrary distribution $`\delta =\left(\mathrm{\Sigma }_k\right)_{kK}`$ in the rest of this paper. We call *processes* the elements of $`K`$. The *location* $`\mathrm{Loc}(a)`$ of an action $`a\mathrm{\Sigma }`$ consists of all processes $`kK`$ such that $`a\mathrm{\Sigma }_k`$: $`\mathrm{Loc}(a)=\{kK|a\mathrm{\Sigma }_k\}`$.
###### Definition 1.2
An *asynchronous automaton* over the distribution $`\left(\mathrm{\Sigma }_k\right)_{kK}`$ consists of a family of finite sets of states $`\left(Q_k\right)_{kK}`$, a family of initial local states $`\left(ı_k\right)_{kK}`$ with $`ı_kQ_k`$, a subset of final global states $`F_{kK}Q_k`$, and a transition relation $`_a_{k\mathrm{Loc}(a)}Q_k\times _{k\mathrm{Loc}(a)}Q_k`$ for each action $`a\mathrm{\Sigma }`$.
The set of *global states* $`Q=_{kK}Q_k`$ can be provided with a set of global transitions $`\stackrel{}{}`$ in such a way that an asynchronous automaton is viewed as a particular automaton. Given an action $`a\mathrm{\Sigma }`$ and two global states $`q=\left(q_k\right)_{kK}`$ and $`r=\left(r_k\right)_{kK}`$, we put $`q\stackrel{a}{}r`$ if $`(\left(q_k\right)_{k\mathrm{Loc}(a)},\left(r_k\right)_{k\mathrm{Loc}(a)})_a`$ and $`q_k=r_k`$ for all $`kK\mathrm{Loc}(a)`$. The initial global state $`ı`$ consists of the collection of initial local states: $`ı=\left(ı_k\right)_{kK}`$. Then the *global automaton* $`𝒜=(Q,ı,\mathrm{\Sigma },\stackrel{}{},F)`$ satisfies Property ID of Def. 1.1. Thus it is an asynchronous system over $`(\mathrm{\Sigma },)`$ and $`L(𝒜)`$ is a regular trace language. An asynchronous automaton is *deterministic* if its global automaton is deterministic, i.e. the local transition relations $`_a`$ are partial functions.
### 1.3 Main result and comparisons to related works
Although deterministic asynchronous automata appear as a restricted subclass of deterministic asynchronous systems, Zielonka’s theorem asserts that any regular trace language can be implemented in the form of a deterministic asynchronous automaton.
###### Theorem 1.3
For any regular trace language $`L`$ there exists a deterministic asynchronous automaton whose global automaton $`𝒜`$ satisfies $`L=L(𝒜)`$.
In a complexity analysis of Zielonka’s construction is detailed. Let $`|Q|`$ be the number of states of the minimal deterministic automaton that accepts $`L`$ and $`|K|`$ be the number of processes. Then the number of local states built by Zielonka’s technique in each process $`kK`$ is $`|Q_k|2^{O(2^{|K|}.|Q|\mathrm{log}|Q|)}`$. The simplified construction by Cori et al. in also suffers from this exponential state-explosion .
Another construction proposed by Pighizzini builds some non-deterministic asynchronous automata from particular rational expressions that refine Ochmański’s theorem . This simpler approach proceeds inductively on the structure of the rational expression. Each step can easily be shown to be polynomial. In particular the number of local states in each process is (at least) *doubled* by each restricted iteration. Consequently in some cases the number of local states in each process is *exponential* in the length of the rational expression.
In the present paper we give a new construction that is *polynomial in $`|Q|`$* (Th. 3.1): It produces $`|Q_k|O(|Q|^d)`$ local states for each process, where $`d=(2.|\mathrm{\Sigma }|+2)^{|\mathrm{\Sigma }|+1}`$, $`|\mathrm{\Sigma }|`$ is the size of $`\mathrm{\Sigma }`$, and $`|Q|`$ is the number of states of some (possibly non-deterministic) asynchronous system that accepts $`L`$. Noteworthy the number of local states $`|Q_k|`$ obtained by our approach is independent from the number of processes $`|K|`$.
## 2 Unfolding algorithm
In the rest of the paper we fix some automaton $`𝒜=(Q,ı,\mathrm{\Sigma },\stackrel{}{},F)`$ that is possibly non-deterministic. The aim of this section is to associate with $`𝒜`$ a family of automata called *boxes* and *triangles* which are defined inductively. The last box built by this construction will be called the *unfolding* of $`𝒜`$ (Def. 2.3).
Boxes and triangles are related to $`𝒜`$ by means of morphisms which are defined as follows. Let $`𝒜_1=(Q_1,ı_1,T,\underset{1}{\overset{}{}},F_1)`$ and $`𝒜_2=(Q_2,ı_2,T,\underset{2}{\overset{}{}},F_2)`$ be two automata over a subset of actions $`T\mathrm{\Sigma }`$. A *morphism $`\sigma :𝒜_1𝒜_2`$ from $`𝒜_1`$ to $`𝒜_2`$* is a mapping $`\sigma :Q_1Q_2`$ from $`Q_1`$ to $`Q_2`$ such that $`\sigma (ı_1)=ı_2`$, $`\sigma (F_1)F_2`$, and $`q_1\underset{1}{\overset{a}{}}q_1^{}`$ implies $`\sigma (q_1)\underset{2}{\overset{a}{}}\sigma (q_1^{})`$. In particular, $`L(𝒜_1)L(𝒜_2)`$.
Now boxes and triangles are associated with an initial state that may not correspond to the initial state of $`𝒜`$. They are associated also with a subset of actions $`T\mathrm{\Sigma }`$. For these reasons, for any state $`qQ`$ and any subset of actions $`T\mathrm{\Sigma }`$, we let $`𝒜_{T,q}`$ denote the automaton $`(Q,q,T,\underset{T}{\overset{}{}},F)`$ where $`\underset{T}{\overset{}{}}`$ is the restriction of $`\stackrel{}{}`$ to the transitions labeled by actions in $`T`$: $`\underset{T}{\overset{}{}}=\stackrel{}{}(Q\times T\times Q)`$.
In this section we shall define the box $`\mathrm{}_{T,q}`$ for all states $`qQ`$ and all subsets of actions $`T\mathrm{\Sigma }`$. The box $`\mathrm{}_{T,q}`$ is a pair $`(_{T,q},\beta _{T,q})`$ where $`_{T,q}`$ is an automaton over $`T`$ and $`\beta _{T,q}:_{T,q}𝒜_{T,q}`$ is a morphism. Similarly, we shall define the triangle $`\mathrm{}_{T,q}`$ for all states $`q`$ and all *non-empty* subsets of transitions $`T`$. The triangle $`\mathrm{}_{T,q}`$ is a pair $`(𝒯_{T,q},\tau _{T,q})`$ where $`𝒯_{T,q}`$ is an automaton over $`T`$ and $`\tau _{T,q}:𝒯_{T,q}𝒜_{T,q}`$ is a morphism.
The *height* of a box $`\mathrm{}_{T,q}`$ or a triangle $`\mathrm{}_{T,q}`$ is the cardinality of $`T`$. Boxes and triangles are defined inductively. We first define the box $`\mathrm{}_{\mathrm{},q}`$ for all states $`qQ`$. Then triangles of height $`h`$ are built upon boxes of height $`g<h`$ and boxes of height $`h`$ are built upon either triangles of height $`h`$ or boxes of height $`g<h`$, whether the dependence graph $`(T,\overline{)})`$ is connected or not.
The base case deals with the boxes of height 0. For all states $`qQ`$, the box $`\mathrm{}_{\mathrm{},q}`$ consists of the morphism $`\beta _{\mathrm{},q}:\{q\}Q`$ that maps $`q`$ to itself together with the automaton $`_{\mathrm{},q}=(\{q\},q,\mathrm{},\mathrm{},F_{\mathrm{},q})`$ where $`F_{\mathrm{},q}=\{q\}`$ if $`qF`$ and $`F_{\mathrm{},q}=\mathrm{}`$ otherwise. More generally a state of a box or a triangle is final if it is associated with a final state of $`𝒜`$.
### 2.1 Building triangles from boxes
Triangles are made of boxes of lower height. Boxes are inserted into a triangle inductively on the height along a tree-like structure and several copies of the same box may appear within a triangle. We want to keep track of this structure in order to prove properties of triangles (and boxes) inductively. This enables us also to allow for the distinction of different copies of the same box within a triangle.
To do this, each state of a triangle is associated with a *rank* $`k`$ such that all states with the same rank come from the same copy of the same box. It is also important to keep track of the height each state comes from, because boxes of a triangle are included inductively on the height. For these reasons, a state of a triangle $`\mathrm{}_{T^{},q^{}}=(𝒯_{T^{},q^{}},\tau _{T^{},q^{}})`$ is encoded as a quadruple $`v=(w,T,q,k)`$ such that $`w`$ is a state from the box $`\mathrm{}_{T,q}`$ with height $`h=|T|`$ and $`v`$ is added to the triangle within the $`k`$-th box inserted into the triangle. Moreover this box is a copy of $`\mathrm{}_{T,q}`$. In that case the state $`v`$ maps to $`\tau _{T^{},q^{}}(v)=\beta _{T,q}(w)`$, that is, the insertion of boxes preserves the correspondance to the states of $`𝒜`$. Moreover the morphism $`\tau _{T^{},q^{}}`$ of a triangle $`\mathrm{}_{T^{},q^{}}`$ is encoded in the data structure of its states.
The construction of the triangle $`\mathrm{}_{T^{},q^{}}`$ is detailed in Algorithm 2.1. It relies on four procedures:
* Build-Box$`(T,q)`$ returns the box $`\mathrm{}_{T,q}`$.
* Mark$`((,\beta ),T,q,k)`$ returns a copy of $`(,\beta )`$ where each state $`w`$ from $``$ is replaced by the marked state $`v=(w,T,q,k)`$.
* Insert$`((𝒯,\tau ),(,\beta ))`$ inserts $`(,\beta )`$ within $`(𝒯,\tau )`$; the initial state of this disjoint union of automata is the initial state of $`(𝒯,\tau )`$.
* Add$`((𝒯,\tau ),(v,a,v^{}))`$ adds a new transition $`v\stackrel{a}{}v^{}`$ to the automaton $`𝒯`$; it is required that $`v`$ and $`v^{}`$ be states of $`𝒯`$.
The construction of the triangle $`\mathrm{}_{T^{},q^{}}`$ starts with building a copy of the base box $`\mathrm{}_{\mathrm{},q^{}}`$ which gets rank $`k=1`$ and whose marked initial state $`(ı_{\mathrm{},\mathrm{},q^{}},\mathrm{},q^{},1)`$ becomes the initial state of $`\mathrm{}_{T^{},q^{}}`$. Along the construction of this triangle, $`k`$ counts the number of boxes already inserted in the triangle. The insertion of boxes proceeds inductively on the height $`h`$ (Line 2.1) as follows: For each state $`v=(w,T,q,l)`$ with height $`|T|=h1`$, if a transition $`\beta _{T,q}(w)\stackrel{a}{}q^{}`$ in $`𝒜`$ carries an action $`aT^{}T`$ (Line 2.1) then a new box $`\mathrm{}_{T^{},q^{}}`$ of height $`h`$ is inserted with $`T^{}=T\{a\}`$ (Line 2.1) and a transition $`v\stackrel{a}{}v^{}`$ is added to the triangle $`𝒯_{T^{},q^{}}`$ in construction (Line 2.1) where $`v^{}`$ is the marked initial state of the new box $`\mathrm{}_{T^{},q^{}}`$. We stress here that $`\tau (v)\stackrel{a}{}\tau (v^{})`$ is a transition of $`𝒜_{T^{},q^{}}`$ because $`\tau (v)=\beta _{T,q}(w)`$ and $`\tau (v^{})=\beta _{T^{},q^{}}(ı_{\mathrm{},T^{},q^{}})=q^{}`$. This observation will show that $`\tau `$ is a morphism. Another useful remark is the following.
###### Lemma 2.1
If a word $`u\mathrm{\Sigma }^{}`$ leads in the triangle $`\mathrm{}_{T^{},q^{}}`$ from its initial state $`(ı_{\mathrm{},\mathrm{},q^{}},\mathrm{},q^{},1)`$ to some state $`v=(w,T,q,l)`$ then each action of $`T`$ occurs in $`u`$.
### 2.2 Building boxes from triangles
We distinguish two cases when we build the box $`\mathrm{}_{T,q}`$ whether the dependence graph $`(T,\overline{)})`$ is connected or not. In case $`(T,\overline{)})`$ is a connected graph then the box $`\mathrm{}_{T,q}`$ collects all triangles $`\mathrm{}_{T,q^{}}`$ for all states $`q^{}Q`$. Each triangle is duplicated a fixed number of times and copies of triangles are connected in some particular way. Similarly to triangles, the states of a box are decorated with a rank $`k`$ that distinguishes states from different triangles and also states from different copies of the same triangle. We adopt the same data structure as for triangles: A state $`v`$ of a box is a quadruple $`(w,T,q,k)`$ where $`w`$ is a state of $`\mathrm{}_{T,q}`$ and $`k`$. Whereas triangles of height $`h`$ are built upon boxes of height $`g<h`$, boxes $`\mathrm{}_{T,q}`$ are built upon triangles $`\mathrm{}_{T,q^{}}`$ with the same set of transitions $`T`$ — and consequently, with the same height. Similarly to the algorithm Build-Triangle, the algorithm that builds boxes uses an integer variable $`k`$ that counts the number of triangles already inserted in the box in construction.
In case the dependence graph $`(T,\overline{)})`$ is not connected, we let $`T_1`$ denote the connected component of $`(T,\overline{)})`$ that contains the least action $`aT`$ w.r.t. the total order $``$ over $`\mathrm{\Sigma }`$ and we put $`T_2=TT_1`$. Then the box $`\mathrm{}_{T,q}`$ is built upon a copy of the box $`\mathrm{}_{T_2,q}`$ connected to copies of boxes $`\mathrm{}_{T_1,q_1}`$ for some states $`q_1Q`$.
The construction of the box $`\mathrm{}_{T^{},q^{}}`$ is detailed in Algorithm 2.1. It relies on ten procedures:
* Base-Box$`(q)`$ returns the base box $`\mathrm{}_{\mathrm{},q}`$.
* Empty-Box returns a special new box called *empty box*.
* Mark, Insert and Add are the procedures used for Build-Triangle. If $`(,\beta )`$ is this special empty box then Insert$`((,\beta ),(𝒯,\tau ))`$ replaces simply $`(,\beta )`$ by $`(𝒯,\tau )`$.
* Missing$`(T^{},q,q^{})`$ returns the set of all pairs $`(w,a)`$ where $`w`$ is a state that has been inserted in the triangle $`\mathrm{}_{T^{},q}`$ within a box $`\mathrm{}_{T^{\prime \prime },q^{\prime \prime }}`$ such that $`|T^{\prime \prime }|=|T^{}|1`$ and the action $`aT^{}T^{\prime \prime }`$ is such that there is a transition $`\tau _{T^{},q}(w)\stackrel{a}{}q^{}`$ in $`𝒜`$ (Alg. 2.1). Due to the structure of triangles, if $`(w,a)`$ is a missing transition then there is no transition $`w\underset{\mathrm{},T^{},q}{\overset{a}{}}w^{}`$ with $`\tau _{T^{},q}(w^{})=q^{}`$ in $`\mathrm{}_{T^{},q}`$.
* Min-Rank$`(T^{},q,,k)`$ returns the minimal rank of a copy of a triangle $`𝒯_{T^{},q}`$ inserted in $``$ where $`k`$ is the maximal rank of triangles in $``$ (Alg. 2.1).
* Max-Out-Degree$`(T^{})`$ returns the number of copies of each triangle $`\mathrm{}_{T^{},q}`$ that should compose the box $`\mathrm{}_{T^{},q^{}}`$. It does not depend on $`q`$ but it depends on the cardinality of all sets Missing$`(T^{},q,q^{})`$ with $`q,q^{}Q`$ (Alg. 2.1). The rôle of these copies is detailed below.
* Clean$`(,\beta )`$ remove all unreachable states from $``$.
* Decomposition$`(T^{})`$ returns the connected component $`T`$ of $`(T^{},\overline{)})`$ that contains the minimal action of $`T^{}`$ w.r.t. the total order $``$.
The construction of the box $`\mathrm{}_{T^{},q^{}}`$ starts with solving the base case where $`T^{}=\mathrm{}`$ (Line 2.1). Assume now that the dependence graph $`(T^{},\overline{)})`$ is connected (Line 2.1). Then the box is initialized as the special empty box (Line 2.1). The number $`m`$ of copies of each triangle $`\mathrm{}_{T^{},q}`$ is computed in Line 2.1 with the help of functions Max-Out-Degree and Missing. Next these copies are inserted and the first copy of $`\mathrm{}_{T^{},q^{}}`$ gets rank $`k=1`$ (Lines 2.1 to 2.1). Consequently the initial state of the box $`\mathrm{}_{T^{},q^{}}`$ in construction is the first copy of the initial state $`ı_{\mathrm{},T^{},q^{}}`$ of the triangle $`\mathrm{}_{T^{},q^{}}`$, that is: $`(ı_{\mathrm{},T^{},q^{}},T^{},q^{},1)`$. Noteworthy copies of the same triangle have consecutive ranks.
In a second step transitions are added to connect these triangles to each other (Lines 2.1 to 2.1). Intuitively a $`a`$-transition is *missing* from the state $`w=(w^{\prime \prime },T^{\prime \prime },q^{\prime \prime },k^{\prime \prime })`$ of the triangle $`\mathrm{}_{T^{},q}`$ to the state $`q^{}`$ of $`𝒜`$ if $`|T^{}T^{\prime \prime }|=1`$ — i.e. this state has been inserted at the highest level in $`\mathrm{}_{T^{},q}`$ — and there exists in $`𝒜`$ a transition $`\tau _{T^{},q}(w)\stackrel{a}{}q^{}`$ with $`aT^{}T^{\prime \prime }`$ but no transition $`w\stackrel{a}{}w^{}`$ with $`\tau _{T^{},q}(w^{})=q^{}`$ in $`\mathrm{}_{T^{},q}`$.
The rôle of Missing is to compute the missing transitions w.r.t. $`q`$, $`q^{}`$, and $`T^{}`$. For each such missing transition $`(w,a)`$ we connect each copy of $`w`$ to the initial state of a copy of $`\mathrm{}_{T^{},q^{}}`$. In this process we require two crucial properties:
* No added transition connects two states from the same copy of the same triangle: $`(w,T^{},q,l)`$ should not be connected to $`(ı_{\mathrm{},T^{},q},T^{},q,l)`$.
* At most one transition connects one copy of $`\mathrm{}_{T^{},q}`$ to one copy of $`\mathrm{}_{T^{},q^{}}`$: If we add from a given copy of $`\mathrm{}_{T^{},q}`$ a transition $`(w_1,T^{},q,l)\stackrel{a}{}(ı_{\mathrm{},T^{},q^{}},T^{},q^{},l^{})`$ and a transition $`(w_2,T^{},q,l)\stackrel{b}{}(ı_{\mathrm{},T^{},q^{}},T^{},q^{},l^{})`$ to the same copy of $`\mathrm{}_{T^{},q^{}}`$ then $`w_1=w_2`$ and $`a=b`$.
The minimal number of copies required to fulfill these conditions is computed by Max-Out-Degree. For a fixed missing transition $`(w,a)`$ from a state $`w`$ of the triangle $`\mathrm{}_{T^{},q}`$ to a state $`q^{}`$ of $`𝒜`$, Lines 2.1 to 2.1 add a transition from the $`j`$-th copy of $`w`$ to the $`c`$-th copy of the initial state of $`\mathrm{}_{T^{},q^{}}`$ with the property that $`jc`$ if $`q=q^{}`$ (Condition $`\text{P}_1`$ above). Moreover states from the $`j`$-th copy of $`\mathrm{}_{T^{},q}`$ are connected to distinct copies of the initial state of $`\mathrm{}_{T^{},q^{}}`$ (Condition $`\text{P}_2`$ above).
Note here that each new transition $`(v,a,v^{})`$ added to $`(,\beta )`$ at Line 2.1 is such that $`\beta (v)\stackrel{a}{}\beta (v^{})`$ is a transition from $`𝒜_{T^{},q^{}}`$ because $`\beta (v)=\tau _{T^{},q}(w)`$, $`\beta (v^{})=\tau _{T^{},q^{}}(ı_{\mathrm{},T^{},q^{}})=q^{}`$, and $`\tau _{T^{},q}(w)\stackrel{a}{}q^{}`$. Again, this observation will show that $`\beta `$ is a morphism. A crucial remark for boxes of connected alphabets is the following.
###### Lemma 2.2
If a non-empty word $`u`$ leads from the initial state of a triangle $`\mathrm{}_{T^{},q}`$ to the initial state of a triangle $`\mathrm{}_{T^{},q^{}}`$ within the box $`\mathrm{}_{T^{},q^{}}`$ then each action of $`T^{}`$ occurs in $`u`$.
For simplicity’s sake our algorithm uses the same number of copies for each triangle. This approach yields in general unreachable states in useless copies. The latter are removed by Clean at Line 2.1.
Assume now that $`(T^{},\overline{)})`$ is not connected (Line 2.1). Let $`T_1`$ be the connected component of $`T^{}`$ that contains the least action of $`T^{}`$ w.r.t. the total order $``$ over $`\mathrm{\Sigma }`$. We put $`T_2=T^{}T_1`$. The construction of the box $`\mathrm{}_{T^{},q^{}}`$ starts with building a copy of the box $`\mathrm{}_{T_2,q^{}}`$. Next for each state $`w`$ of $`\mathrm{}_{T_2,q^{}}`$ and each transition $`\beta _{T_2,q}(w)\stackrel{a}{}q^{}`$ with $`aT_1`$, the algorithm inserts a (new) copy of the box $`\mathrm{}_{T_1,q^{}}`$ and adds a transition from the copy of $`w`$ to the initial state of the copy of $`\mathrm{}_{T_1,q^{}}`$. By recursive calls of Build-Box the box $`\mathrm{}_{T^{},q}`$ is built along a tree-like structure upon copies of boxes $`\mathrm{}_{T^{},q^{}}`$ where $`T^{}`$ is a connected component of $`T^{}`$.
### 2.3 Remarks
From a mathematical viewpoint, Algorithms 2.1 to 2.1 are meant to define boxes $`\mathrm{}_{T,q}`$ and triangles $`\mathrm{}_{T,q}`$. Thus two instances of Build-Triangle$`(T,q)`$ produce the same object. For this reason, we speak of *the* triangle $`\mathrm{}_{T,q}`$. This is particularly important to understand the interaction between Build-Box and Missing. In case $`T`$ is connected, Algorithm Build-Box proceeds in two steps. First several copies of each triangle $`\mathrm{}_{T,q}`$ are collected and next some transitions are added from some states of copies of $`\mathrm{}_{T,q}`$ to the initial state of copies of $`\mathrm{}_{T,q^{}}`$. These additional transitions are computed in a separate function Missing that depends on triangles. It is crucial that the triangles $`\mathrm{}_{T,q}`$ used by the function Missing be the same as the triangles $`\mathrm{}_{T,q}`$ inserted in Build-Box.
From a more computational viewpoint, Algorithms 2.1 to 2.1 can obviously be implemented. To do this, we require that each triangle and each box be constructed only once. An alternative to this requirement is to adapt the parameters of the function Missing and ensure that Build-Box transfers its own triangle $`\mathrm{}_{T,q}`$ instead of the pair $`(T,q)`$ to that function so that the set of states computed by Missing matches the set of states used by Build-Box. However it need not to transfert its own triangle $`\mathrm{}_{T,q}`$ to the function Max-Out-Degree because this function works on triangles up to isomorphisms.
In this section we have built a family of boxes and triangles from a fixed automaton $`𝒜`$. This construction leads us to the definition of the unfolding of $`𝒜`$ as follows.
###### Definition 2.3
The *unfolding* $`𝒜_{\mathrm{Unf}}`$ of the automaton $`𝒜=(Q,ı,\mathrm{\Sigma },\stackrel{}{},F)`$ is the box $`_{\mathrm{\Sigma },ı}`$; moreover $`\beta _{\mathrm{Unf}}`$ denote the mapping $`\beta _{\mathrm{\Sigma },ı}`$ from the states of $`𝒜_{\mathrm{Unf}}`$ to $`Q`$.
In the next section we study some complexity, structural, and semantical properties of this object. We assume that $`𝒜`$ satisfies Property ID of Definition 1.1 so that it accepts a regular trace language $`L`$. We explain how to build from the unfolding $`𝒜_{\mathrm{Unf}}`$ a non-deterministic asynchronous automaton that accepts $`L(𝒜)`$.
## 3 Properties of the unfolding algorithm
In this section we fix a regular trace language $`L`$ over the independence alphabet $`(\mathrm{\Sigma },)`$. We assume that the possibly non-deterministic automaton $`𝒜`$ fulfills Property ID of Def. 1.1 and satisfies $`L(𝒜)=L`$. First we sketch a complexity analysis of the number of states in the unfolding $`𝒜_{\mathrm{Unf}}`$. Next we show in Subsection 3.2 how to build from $`𝒜_{\mathrm{Unf}}`$ an asynchronous automaton whose global automaton accepts $`L(𝒜)`$.
### 3.1 Complexity analysis
For all naturals $`n0`$ we denote by $`\beta _n`$ the maximal number of states in a box $`_{T,q}`$ with $`|T|=n`$ and $`qQ`$. Similarly for all naturals $`n1`$ we denote by $`\tau _n`$ the maximal number of states in a triangle $`𝒯_{T,q}`$ with $`|T|=n`$ and $`qQ`$. Noteworthy $`\beta _0=1`$ and $`\tau _1=1`$. Moreover $`\tau _n`$ is non-decreasing because the triangle $`\mathrm{}_{T^{},q}`$ is a subautomaton of the triangle $`\mathrm{}_{T,q}`$ as soon as $`T^{}T`$. In the following we assume $`2n|\mathrm{\Sigma }|`$.
Consider some subset $`T\mathrm{\Sigma }`$ with $`|T|=n`$. Each triangle $`𝒯_{T,q}`$ is built inductively upon boxes of height $`hn1`$ (see Alg. 2.1). We distinguish two cases. First boxes of height $`h<n1`$ are inserted. Each of these boxes appears also in some triangle $`𝒯_{T^{},q}`$ with $`T^{}T`$ and $`|T^{}|=n1`$. Each of these triangles is a subautomaton of $`𝒯_{T,q}`$ with at most $`\tau _{n1}`$ states. Moreover there are only $`n`$ such triangles which give rise to at most $`n.\tau _{n1}`$ states built along this first step. Second, boxes of height $`n1`$ are inserted and connected to states inserted at height $`n2`$. Each of these states belongs to some box $`\mathrm{}_{T^{},q^{}}`$ with $`|T^{}|=n2`$; it gives rise to at most $`2.|Q|`$ boxes at height $`n1`$ because $`|TT^{}|=2`$: This yields at most $`2.|Q|.\beta _{n1}`$ new states. Altogether we get
$$\tau _nn.\tau _{n1}.(1+2.|Q|.\beta _{n1})|\mathrm{\Sigma }|.\tau _{n1}.3.|Q|.\beta _{n1}$$
(1)
Consider now a connected subset $`T\mathrm{\Sigma }`$ with $`|T|=n1`$. Then each box $`_{T,q}`$ is built upon triangles $`𝒯_{T,q^{}}`$ of height $`n1`$ (see Alg. 2.1). We can check that the value $`m=`$Max-Out-Degree$`(T)`$ is at most $`\tau _{n1}+1`$. Therefore the box $`_{T,q}`$ contains at most $`2.\tau _{n1}`$ copies of each triangle $`𝒯_{T,q^{}}`$. Hence $`\beta _{n1}2.|Q|.\tau _{n1}^2`$.
Consider now a non-connected subset $`T\mathrm{\Sigma }`$ with $`|T|=n1`$. Then each box $`_{T,q}`$ is built upon copies of boxes $`_{T^{},q^{}}`$ where $`T^{}`$ is a connected component of $`(T,\overline{)})`$. These boxes are inserted inductively along recursive calls of Build-Box and they are connected in a tree-like structure. Each of these boxes contains at most $`2.|Q|.\tau _{n2}^2`$ states as explained above. From each state of these boxes at most $`(n2).|Q|`$ new boxes are connected. Thus each box $`_{T^{},q^{}}`$ is connected to at most $`c=|\mathrm{\Sigma }|.2.|Q|^2.\tau _{n2}^2`$ boxes in the tree-like structure. Consequently there are at most $`1+c+c^2+c^3+\mathrm{}+c^{n2}`$ boxes. It follows that
$$\beta _{n1}c^{n1}.2.|Q|.\tau _{n2}^22^n.|\mathrm{\Sigma }|^{n1}.|Q|^{2.n1}.\tau _{n2}^{2.n}$$
(2)
Since $`\tau _{n2}\tau _{n1}`$ we get in both cases $`\beta _{n1}2^{|\mathrm{\Sigma }|}.|\mathrm{\Sigma }|^{|\mathrm{\Sigma }|1}.|Q|^{2.|\mathrm{\Sigma }|1}.(\tau _{n1})^{2.|\mathrm{\Sigma }|}`$.
We can now apply $`(1)`$ and get $`\tau _nN.\tau _{n1}^d`$ where $`N=3.2^{|\mathrm{\Sigma }|}.|\mathrm{\Sigma }|^{|\mathrm{\Sigma }|}.|Q|^{2.|\mathrm{\Sigma }|}`$ and $`d=2.|\mathrm{\Sigma }|+1`$. Since $`\tau _1=1`$, we get $`\tau _nN^{d^{n1}}`$. We can apply $`(2)`$ with $`n`$ instead of $`n1`$ and get $`\beta _n2.|Q|.N.(\tau _n)^{2.(n+1)}2.|Q|.N.N^{d^{n1}.(2n+2)}`$. Finally we have
$$\beta _{|\mathrm{\Sigma }|}2.|Q|.(3.2^{|\mathrm{\Sigma }|}.|\mathrm{\Sigma }|^{|\mathrm{\Sigma }|}.|Q|^{2.|\mathrm{\Sigma }|})^{(2.|\mathrm{\Sigma }|+2)^{|\mathrm{\Sigma }|}}O(|Q|^{(2.|\mathrm{\Sigma }|+2)^{|\mathrm{\Sigma }|+1}})$$
(3)
### 3.2 Construction of an asynchronous automaton
Finally we build from the unfolding $`𝒜_{\mathrm{Unf}}=(Q_{\mathrm{Unf}},ı_{\mathrm{Unf}},\mathrm{\Sigma },\underset{\mathrm{Unf}}{\overset{}{}},F_{\mathrm{Unf}})`$ of $`𝒜`$ an asynchronous automaton $`\widehat{𝒜_{\mathrm{Unf}}}`$ that accepts $`L(𝒜)`$. We define $`\widehat{𝒜_{\mathrm{Unf}}}`$ as follows. First we put $`Q_k=Q_{\mathrm{Unf}}`$ for each process $`kK`$. Next the initial state is the $`|K|`$-tuple $`(ı_{\mathrm{Unf}},\mathrm{},ı_{\mathrm{Unf}})`$. Moreover for each action $`a`$, the pair $`(\left(q_k\right)_{k\mathrm{Loc}(a)},\left(r_k\right)_{k\mathrm{Loc}(a)})`$ belongs to the transition relation $`_a`$ if there exist two states $`q,rQ_{\mathrm{Unf}}`$ and a transition $`q\underset{\mathrm{Unf}}{\overset{a}{}}r`$ in the unfolding such that the two following conditions are satisfied:
* for all $`k\mathrm{Loc}(a)`$, $`q_k\underset{\mathrm{Unf}}{\overset{u}{}}q`$ for some word $`u(\mathrm{\Sigma }\mathrm{\Sigma }_k)^{}`$;
* for all $`k\mathrm{Loc}(a)`$, $`r_k=r`$; in particular all $`r_k`$ are equal.
Finally, a global state $`\left(q_k\right)_{kK}`$ is *final* if there exists a final state $`qQ_{\mathrm{Unf}}`$ such that for all $`kK`$ there exists a path $`q_k\underset{\mathrm{Unf}}{\overset{u}{}}q`$ for some word $`u(\mathrm{\Sigma }\mathrm{\Sigma }_k)^{}`$.
###### Theorem 3.1
The asynchronous automaton $`\widehat{𝒜_{\mathrm{Unf}}}`$ satisfies $`L(\widehat{𝒜_{\mathrm{Unf}}})=L(𝒜)`$. Moreover the number of local states $`Q_k`$ in each process is polynomial in $`|Q|`$ where $`|Q|`$ is the number of states in $`𝒜`$; more precisely $`|Q_k|O(|Q|^d)`$ where $`d=(2.|\mathrm{\Sigma }|+2)^{|\mathrm{\Sigma }|+1}`$.
### 3.3 Sketch of proof
By induction on the structure of the unfolding it is not difficult to check the following first property (see Appendix 0.A).
###### Lemma 3.2
The mapping $`\beta _{\mathrm{Unf}}`$ is a morphism from the unfolding $`𝒜_{\mathrm{Unf}}`$ to $`𝒜`$. Moreover for all $`uL(𝒜)`$ there exists $`vL(𝒜_{\mathrm{Unf}})`$ such that $`vu`$.
The proof of Theorem 3.1 relies on an intermediate asynchronous automaton $`\overline{𝒜_{\mathrm{Unf}}}`$ over some extended independence alphabet $`(\overline{\mathrm{\Sigma }},\overline{})`$. We consider the alphabet $`\overline{\mathrm{\Sigma }}=\mathrm{\Sigma }\{(a,k)\mathrm{\Sigma }\times K|a\mathrm{\Sigma }_k\}`$ provided with the independence relation $`\overline{}`$ such that $`a\overline{}\overline{)}b`$ iff $`a\overline{)}b`$, $`a\overline{}\overline{)}(b,k)`$ iff $`a\mathrm{\Sigma }_k`$, and $`(a,k)\overline{}\overline{)}(b,k^{})`$ iff $`k=k^{}`$ for all actions $`a,b\mathrm{\Sigma }`$ and all processes $`k,k^{}K`$. For each process $`kK`$ we put $`\overline{\mathrm{\Sigma }_k}=\mathrm{\Sigma }_k\{(a,k)|a\mathrm{\Sigma }\mathrm{\Sigma }_k\}`$. It is easy to check that $`\left(\overline{\mathrm{\Sigma }}_k\right)_{kK}`$ is a distribution of $`(\overline{\mathrm{\Sigma }},\overline{})`$. Now $`\overline{𝒜_{\mathrm{Unf}}}`$ shares with $`\widehat{𝒜_{\mathrm{Unf}}}`$ its local states $`Q_k`$ and its initial state. A global state $`\left(q_k\right)_{kK}`$ is final if there exists a final state $`qF_{\mathrm{Unf}}`$ of the unfolding such that $`q_k=q`$ for all $`kK`$. For each action $`a\mathrm{\Sigma }`$ its transition relation $`\overline{}_a`$ is such that $`((q_k)_{k\mathrm{Loc}(a)},(q_k^{})_{k\mathrm{Loc}(a)})\overline{}_a`$ if there exists a transition $`q\underset{\mathrm{Unf}}{\overset{a}{}}q^{}`$ such that $`q_k=q`$ and $`q_k^{}=q^{}`$ for all $`k\mathrm{Loc}(a)`$. Moreover for each internal action $`(a,k)\overline{\mathrm{\Sigma }}\mathrm{\Sigma }`$, we put $`(q,q^{})\overline{}_{(a,k)}`$ if $`q\underset{\mathrm{Unf}}{\overset{a}{}}q^{}`$.
We consider the projection morphism $`\rho :\overline{\mathrm{\Sigma }}^{}\mathrm{\Sigma }^{}`$ such that $`\rho (\epsilon )=\epsilon `$, $`\rho (u.a)=\rho (u).a`$ if $`a\mathrm{\Sigma }`$, and $`\rho (u.a)=\rho (u)`$ if $`a\overline{\mathrm{\Sigma }}\mathrm{\Sigma }`$. It is not difficult to prove that $`L(𝒜_{\mathrm{Unf}})\rho (L(\overline{𝒜_{\mathrm{Unf}}}))`$ and $`\rho (L(\overline{𝒜_{\mathrm{Unf}}}))=L(\widehat{𝒜_{\mathrm{Unf}}})`$. These two basic properties do not rely on the particular structure of $`𝒜_{\mathrm{Unf}}`$: They hold actually for any automaton.
On the contrary the proof of the next lemma is very technical and tedious and relies on the particular construction of boxes and triangles (see Appendix 0.B for some details).
###### Lemma 3.3
For each box $`\mathrm{}_{T,q}=(_{T,q},\beta _{T,q})`$ we have $`\rho (L(\overline{_{T,q}}))[L(_{T,q})]`$.
We can now conclude. By Lemma 3.2 we have $`[L(𝒜_{\mathrm{Unf}})]=L(𝒜)`$. On the other hand the two basic properties show that $`[L(𝒜_{\mathrm{Unf}})]L(\widehat{𝒜_{\mathrm{Unf}}})`$. Conversely Lemma 3.3 yields $`L(\widehat{𝒜_{\mathrm{Unf}}})=\rho (L(\overline{𝒜_{\mathrm{Unf}}}))[L(𝒜_{\mathrm{Unf}})]`$. Therefore $`L(𝒜)=L(\widehat{𝒜_{\mathrm{Unf}}})`$.
## Conclusion and future work
We have presented a polynomial algorithm for the construction of non-deterministic asynchronous automata from regular trace languages. We have shown that this new unfolding method improves the complexity of known techniques in terms of the number of local and global states. Several variations of our approach lead to analogous complexity results. We have selected here the simplest version to analyse. But it might not be the more efficient in practice.
Interestingly this unfolding method can be adapted to the implementation of any globally-cooperative compositional high-level message sequence charts as investigated in . At present we are developping more involved unfolding techniques in order to construct deterministic safe asynchronous automata . We are also investigating a possible extension of our unfolding technique to infinite traces .
## Appendix 0.A Proof of Lemma 3.2
An immediate induction shows that for each box $`\mathrm{}_{T,q}=(_{T,q},\beta _{T,q})`$ and each triangle $`\mathrm{}_{T,q}=(𝒯_{T,q},\tau _{T,q})`$, the mappings $`\beta _{T,q}`$ and $`\tau _{T,q}`$ are morphisms from $`_{T,q}`$ to $`𝒜_{T,q}`$ and from $`𝒯_{T,q}`$ to $`𝒜_{T,q}`$ respectively. In particular $`L(𝒜_{\mathrm{Unf}})L(𝒜)`$.
By induction on the size of $`T`$ we prove that for all paths $`q\stackrel{u}{}q_1`$ of $`𝒜_{T,q}`$ there exists an equivalent word $`u^{}u`$ such that $`ı_{\mathrm{},T,q}\stackrel{u^{}}{}v`$ is a path of $`_{T,q}`$ and $`\beta _{T,q}(v)=q_1`$. This property is trivial for the empty set $`T=\mathrm{}`$ because $`𝒜_{\mathrm{},q}`$ and $`_{\mathrm{},q}`$ are reduced to the state $`q`$. We shall distinguish two cases whether $`T`$ is connected or not.
Assume first that $`T`$ is a connected set of actions. We proceed by induction on the length of $`u`$. The property holds for the empty word because $`\beta _{T,q}(ı_{\mathrm{},T,q})=q`$. Let $`q\stackrel{u}{}q_1\stackrel{a}{}q_2`$ be a path of $`𝒜_{T,q}`$. By induction there is $`u^{}u`$ such that $`ı_{\mathrm{},T,q}\stackrel{u^{}}{}v`$ is a path in $`_{T,q}`$ and $`\beta _{T,q}(v)=q_1`$. Then, by construction, the state $`v=(w,T,q^{},k)`$ comes from some triangle $`\mathrm{}_{T,q^{}}`$. Furthermore $`w`$ comes from a box $`\mathrm{}_{T^{\prime \prime },q^{\prime \prime }}`$ inserted in $`\mathrm{}_{T,q^{}}`$: We have $`w=(w^{\prime \prime },T^{\prime \prime },q^{\prime \prime },k^{\prime \prime })`$. We distinguish several cases.
1. If $`aTT^{\prime \prime }`$ and $`|TT^{\prime \prime }|=1`$. Then $`(w,a)`$ belongs to $`\text{Missing}(T,q^{},q_2)`$. Consequently Line $`25`$ of Alg. 2.1 shows that $`v\stackrel{a}{}(ı_{\mathrm{},T,q_2},T,q_2,k^{})`$ for some integer $`k^{}`$ and $`\beta _{T,q}(ı_{\mathrm{},T,q_2},T,q_2,k^{})=\tau _{T,q_2}(ı_{\mathrm{},T,q_2})=q_2`$.
2. If $`aTT^{\prime \prime }`$ and $`|TT^{\prime \prime }|1`$. Then Line $`13`$ of Alg. 2.1 shows that $`w\stackrel{a}{}w^{}`$ with $`w^{}=(ı_{\mathrm{},T^{\prime \prime }\{a\},q_2},T^{\prime \prime }\{a\},q_2,k^{})`$ for some integer $`k^{}`$ is a transition of $`𝒯_{T,q^{}}`$ and $`\tau _{T,q^{}}(w^{})=q_2`$. Consequently $`v\stackrel{a}{}(w^{},T,q^{},k)`$ is a transition of $`_{T,q}`$ and $`\beta _{T,q}(w^{},T,q^{},k)=q_2`$ (see Line $`14`$ of Alg. 2.1).
3. If $`aT^{\prime \prime }`$. By construction the path $`ı_{\mathrm{},T,q}\stackrel{u^{}}{}v`$ of $`_{T,q}`$ consists of the sequence of transitions $`ı_{\mathrm{},T,q}\stackrel{u_1}{}v_1\stackrel{u_2}{}v`$ such that $`u_1.u_2=u^{}`$, $`v_1=(w_1,T,q^{},k)`$, $`w_1=(ı_{\mathrm{},T^{\prime \prime },q^{\prime \prime }},T^{\prime \prime },q^{\prime \prime },k^{\prime \prime })`$ and all states $`v_2`$ reach along the path $`v_1\stackrel{u_2}{}v`$ come from the same box $`\mathrm{}_{T^{\prime \prime },q^{\prime \prime }}`$ of the same triangle $`\mathrm{}_{T,q^{}}`$ that is $`v_2`$ is some tuple $`((w_2,T^{\prime \prime },q^{\prime \prime },k^{\prime \prime }),T,q^{},k)`$. Consequently each action $`b`$ that occurs in $`u_2`$ belongs to $`T^{\prime \prime }`$: It follows that $`q^{\prime \prime }\stackrel{u_2}{}q_1\stackrel{a}{}q_2`$ is a path of $`𝒜_{T^{\prime \prime },q^{\prime \prime }}`$. By induction there is an equivalent word $`u_2^{}u_2.a`$ such that $`ı_{\mathrm{},T^{\prime \prime },q^{\prime \prime }}\stackrel{u_2^{}}{}w^{}`$ is a path of $`_{T^{\prime \prime },q^{\prime \prime }}`$ and $`\beta _{T^{\prime \prime },q^{\prime \prime }}(w^{})=q_2`$. Consequently $`v_1\stackrel{u_2^{}}{}((w^{},T^{\prime \prime },q^{\prime \prime },k^{\prime \prime }),T,q^{},k)`$ is a path of $`_{T,q}`$ and $`\beta _{T,q}((w^{},T^{\prime \prime },q^{\prime \prime },k^{\prime \prime }),T,q^{},k)=q_2`$ (see Line $`12`$ of Alg. 2.1 and Line $`14`$ of Alg. 2.1).
Suppose now that $`T`$ is an unconnected set of actions. Let $`q\stackrel{u}{}q_1`$ be a path of $`𝒜_{T,q}`$, $`T_1`$ be the connected component that contains the least action of $`T`$ and $`T_2=TT_1`$. If $`u|T_1=\epsilon `$ then $`q\stackrel{u}{}q_1`$ is also a path of $`𝒜_{T_2,q}`$. Consequently, by induction there exists $`u^{}u`$ such that $`ı_{\mathrm{},T_2,q}\stackrel{u^{}}{}w`$ is a path of $`_{T_2,q}`$ and $`\beta _{T_2,q}(w)=q_1`$. It follows by Line $`33`$ of Alg. 2.1 that $`ı_{\mathrm{},T,q}\stackrel{u^{}}{}(w,T_2,q,1)`$ is a path of $`_{T,q}`$ and $`\beta _{T,q}(w,T_2,q,1)=q_1`$. If $`u|T_1=a.u_1`$ and $`u|T_2=u_2`$ then $`q\stackrel{u_2}{}q_2\stackrel{a}{}q_3\stackrel{u_1}{}q_1`$ is also a path of $`𝒜_{T,q}`$ because $`u_2.a.u_1u`$ and $`𝒜`$ satisfies ID. Moreover $`q\stackrel{u_2}{}q_2`$ is a path of $`𝒜_{T_2,q}`$ and $`q_3\stackrel{u_1}{}q_1`$ is a path of $`𝒜_{T_1,q_3}`$. Consequently, by induction, there exists $`u_2^{}u_2`$ such that $`ı_{\mathrm{},T_2,q}\stackrel{u_2^{}}{}w_2`$ is a path of $`_{T_2,q}`$ and $`\beta _{T_2,q}(w_2)=q_2`$, and on other hand, there exists $`u_1^{}u_1`$ such that $`ı_{\mathrm{},T_1,q_3}\stackrel{u_1^{}}{}w_1`$ is a path of $`_{T_1,q_3}`$ and $`\beta _{T_1,q_3}(w_1)=q_1`$. Then Alg. 2.1 ensures that $`ı_{\mathrm{},T,q}\stackrel{u_2^{}}{}(w_2,T_2,q,1)`$ is a path of $`_{T,q}`$, $`\beta _{T,q}(w_2,T_2,q,1)=q_2`$ (Line $`33`$), $`(ı_{\mathrm{},T_1,q_3},T_1,q_3,k)\stackrel{u_1^{}}{}(w_1,T_1,q_3,k)`$ is a path of $`_{T,q}`$ for some integer $`k`$ and $`\beta _{T,q}(w_1,T_1,q_3,k)=q_1`$ (Line $`40`$). Finally we have $`(w_2,T_2,q,1)\stackrel{a}{}(ı_{\mathrm{},T_1,q_3},T_1,q_3,k)`$ (Line $`35`$ and $`41`$). It follows that $`ı_{\mathrm{},T,q}\stackrel{u_2^{}.a.u_1^{}}{}(w_1,T_1,q_3,k)`$ is a path of $`_{T,q}`$, $`\beta _{T,q}(w_1,T_1,q_3,k)=q_1`$ and $`u_2^{}.a.u_1^{}u`$.
## Appendix 0.B Proof sketch of Lemma 3.3
The complete proof of Lemma 3.3 requires about 20 pages of tedious technical details. In this appendix we present the two main ideas that lead the argument. We need first to introduce some basic definitions and notations precisely.
#### Some basic definitions and notations.
Let $`𝒜`$ be some automaton over $`\mathrm{\Sigma }`$. A *path* of length $`n\{0\}`$ is a sequence of transitions $`\left(q_i\stackrel{a_i}{}q_i^{}\right)_{i[1,n]}`$ such that $`q_i^{}=q_{i+1}`$ for all integers $`0<i<n`$. For all words $`u\mathrm{\Sigma }^{}`$ we write $`q\stackrel{u}{}q^{}`$ to denote a path $`\left(q_i\stackrel{a_i}{}q_i^{}\right)_{i[1,n]}`$ where $`q_1=q`$, $`q_n^{}=q^{}`$, and $`u=a_1\mathrm{}a_n`$. Then $`q`$ is called the domain of $`q\stackrel{u}{}q^{}`$ and $`q^{}`$ is called its codomain. A path of length 0 is simply a state $`q`$ of $`𝒜`$. Its domain and codomain are equal to $`q`$.
If $`s`$ and $`s^{}`$ are two paths such that the codomain of $`s`$ is the domain of $`s^{}`$ then the *product* $`ss^{}`$ is defined in a natural way: If the length of $`s`$ is 0 then $`ss^{}=s^{}`$; if the length of $`s^{}`$ is 0 then $`ss^{}=s`$; otherwise $`ss^{}`$ is the concatenation of $`s`$ and $`s^{}`$.
Note that if $`s`$ is a path of the length $`l>0`$ then it is the product of two paths $`s=s_1s_2`$ where the length of $`s_1`$ is $`1`$ and the length of $`s_2`$ is $`l1`$. Moreover such a product is unique. This remark allows us to define mappings for paths inductively on the length.
#### Projections of global states and executions.
Assume now that $`𝒜`$ is (the global system of) an asynchronous automaton over the distribution $`\left(\mathrm{\Sigma }_k\right)_{kK}`$. Then a path of $`𝒜`$ is called an *execution*. For convenience we shall consider the component automata $`\left(𝒜_k\right)_{kK}`$ defined as follows: For each process $`jK`$, $`𝒜_j=(Q_j,ı_j,\mathrm{\Sigma }_j,\underset{j}{\overset{}{}},Q_j)`$ where $`q_j\underset{j}{\overset{a}{}}q_j^{}`$ if there are $`q=\left(q_k\right)_{k\mathrm{Loc}(a)}`$ and $`q=\left(q_k^{}\right)_{k\mathrm{Loc}(a)}`$ such that $`(q,q^{})_a`$. Note here that $`j\mathrm{Loc}(a)`$ since $`a\mathrm{\Sigma }_j`$.
Now the *projection $`s|k`$* of an execution $`s`$ of $`𝒜`$ onto a process $`jK`$ is a path of $`𝒜_j`$ defined inductively as follows:
* $`s|j=q_j`$ if $`s`$ is a path of length 0 that corresponds to the global state $`(q_k)_{kK}`$;
* $`s|j=q_j\stackrel{a}{}q_j^{}(s^{}|j)`$ if $`s`$ is the product $`s=ts^{}`$ where $`t`$ is a transition $`(q_k)_{kK}\stackrel{a}{}(q_k^{})_{kK}`$ and $`j\mathrm{Loc}(a)`$.
* $`s|j=s^{}|j`$ if $`s`$ is the product $`s=ts^{}`$ where $`t`$ is a transition $`q\stackrel{a}{}q^{}`$ and $`j\mathrm{Loc}(a)`$.
#### Executions of extended asynchronous automata.
In the paper we define the *extended asynchronous automaton* $`\overline{𝒜_{\mathrm{Unf}}}`$ of the unfolding automaton $`𝒜_{\mathrm{Unf}}`$. This definition can naturally be generalized to any automaton. Let $`𝒜`$ be an automaton and $`\overline{𝒜}`$ the corresponding extended asynchronous automaton. We say that an execution $`s=q\stackrel{u}{}q^{}`$ of $`\overline{𝒜}`$ is *arched* if there are two states $`v`$ and $`v^{}`$ in $`𝒜`$ such that for all $`kK`$, $`q|k=v`$ and $`q^{}|k=v^{}`$. Noteworthy each execution that leads an extended asynchronous automaton $`\overline{𝒜}`$ from its global initial state to some global final state is arched.
We define now a function $`\gamma `$ that associates each action of $`\overline{\mathrm{\Sigma }}`$ to the corresponding action of $`\mathrm{\Sigma }`$ in a natural way: For all actions $`(a,k)\overline{\mathrm{\Sigma }}\mathrm{\Sigma }`$, $`\gamma (a,k)=a`$ and for all actions $`a\mathrm{\Sigma }`$, $`\gamma (a)=a`$. As usual this map extends from actions to words and we get $`\gamma :\overline{\mathrm{\Sigma }}^{}\mathrm{\Sigma }^{}`$. We can also extend the mapping $`\gamma `$ as a function from paths of component automata $`𝒜_k`$ to paths of $`𝒜`$ as follows. For each sequence $`s`$ that is a path of some $`𝒜_k`$, we define $`\gamma (s)`$ inductively on the length of $`s`$ by
* $`\gamma (s)=q`$ if the length of $`s`$ is 0 and $`s=q`$.
* $`\gamma (s)=q\stackrel{\gamma (a)}{}q^{}\gamma (s^{})`$ if $`s`$ is a product $`s=ts^{}`$ where $`t`$ is a transition $`q\stackrel{a}{}q^{}`$.
Clearly if $`s`$ is an execution of $`\overline{𝒜}`$ and $`k`$ a process of $`K`$ then $`s|k`$ is a path of $`𝒜_k`$ and $`\gamma (s|k)`$ is a path of $`𝒜`$.
#### Definitions associated to unfoldings.
Let $`T`$ be a non-empty subset of $`\mathrm{\Sigma }`$. We consider the triangle $`𝒯_{T,q}=(Q_{\mathrm{},T,q},ı_{\mathrm{},T,q},\underset{\mathrm{},T,q}{\overset{}{}},F_{\mathrm{},T,q})`$. Let $`v`$ be a state from $`𝒯_{T,q}`$. By construction of $`𝒯_{T,q}`$, $`v`$ is a quadruple $`(w,T^{},q^{},k^{})`$ such that $`w`$ is a state from the box $`\mathrm{}_{T^{},q^{}}`$ and $`k^{}`$. We say that the *box location* of $`v`$ is $`l^{\mathrm{}}(v)=(T^{},q^{},k^{})`$. We define the *sequence of boxes* reached along a path $`s=q\stackrel{u}{}q^{}`$ in $`𝒯_{T,q}`$ as follows:
* If the length of $`s`$ is $`0`$ and $`s`$ corresponds to state $`qQ_{\mathrm{},T,q}`$ then $`𝕃^{\mathrm{}}(s)=l^{\mathrm{}}(q)`$.
* If $`s`$ is a product $`s=s^{}t`$ where $`t`$ is the transition $`q\stackrel{a}{}q^{}`$ then two cases appear:
+ If $`l^{\mathrm{}}(q)=l^{\mathrm{}}(q^{})`$ then $`𝕃^{\mathrm{}}(s)=𝕃^{\mathrm{}}(s^{})`$;
+ If $`l^{\mathrm{}}(q)l^{\mathrm{}}(q^{})`$ then $`𝕃^{\mathrm{}}(s)=𝕃^{\mathrm{}}(s^{}).l^{\mathrm{}}(q^{})`$
Similarly we define the sequence of triangles $`𝕃^{\mathrm{}}(s_1)`$ reached by a path $`s_1`$ in a box $`_{T_1,q_1}`$ where $`T_1`$ is a non-empty *connected* set of actions and the sequence of boxes $`𝕃^{\mathrm{}}(s_2)`$ reached by a path $`s_2`$ in a box $`_{T_2,q_2}`$ where $`T_2`$ is an *unconnected* set of actions.
#### Two main properties of unfoldings.
The following proposition states that all processes behave similarly in an extended asynchronous automaton built from boxes or triangles.
###### Proposition 0.B.1
Let $`_{T_1,q_1}`$ be a box with $`T_1`$ a non-empty *connected* set of actions, $`_{T_2,q_2}`$ be a box with $`T_2`$ an *unconnected* set of actions, and $`𝒯_{T_3,q_3}`$ be a triangle with $`T_3`$ a non-empty set of actions. Let $`s_1`$, $`s_2`$ and $`s_3`$ be *arched* executions of $`\overline{}_{T_1,q_1}`$, $`\overline{}_{T_2,q_2}`$ and $`\overline{𝒯}_{T_3,q_3}`$ respectively. Then:
1. $`k,k^{}\mathrm{Loc}(T_1)`$, $`𝕃^{\mathrm{}}(\gamma (s_1|k))=𝕃^{\mathrm{}}(\gamma (s_1|k^{}))`$;
2. $`k,k^{}K`$, $`𝕃^{\mathrm{}}(\gamma (s_2|k))=𝕃^{\mathrm{}}(\gamma (s_2|k^{}))`$;
3. $`k,k^{}K`$, $`𝕃^{\mathrm{}}(\gamma (s_3|k))=𝕃^{\mathrm{}}(\gamma (s_3|k^{}))`$.
Proof. Property 2 and Property 3 stem from the remark that $`_{T_2,q_2}`$ and $`𝒯_{T_3,q_3}`$ are made of boxes connected along a tree-like structure. The proof of Property 1 is more subtle. Let $`a`$ be an action of $`T_1`$ and $`k,k^{}`$ be two processes of $`\mathrm{Loc}(a)`$. We proceed by contradiction. Let $`𝒯`$ and $`𝒯^{}`$ be the first triangles that differ in $`𝕃^{\mathrm{}}(\gamma (s_1|k))`$ and $`𝕃^{\mathrm{}}(\gamma (s_1|k^{}))`$. Let $`c`$ be the number of $`a`$-transitions that occur in $`s_1`$ just before $`\gamma (s_1|k)`$ and $`\gamma (s_1|k^{})`$ reach $`𝒯`$ and $`𝒯^{}`$. Since $`s_1`$ is arched, $`\gamma (s_1|k)`$ and $`\gamma (s_1|k^{})`$ have to meet eventually for the last state. Therefore $`\gamma (s_1|k)`$ and $`\gamma (s_1|k^{})`$ have to leave triangles $`𝒯`$ and $`𝒯^{}`$ respectively. Consequently, there is a $`(c+1)^{th}`$ $`a`$-transition $`q\stackrel{a}{}q^{}`$ in $`s_1`$. Moreover, this transition is such that $`q|k`$ is a state from $`𝒯`$ whereas $`q|k^{}`$ is a state from $`𝒯^{}`$, that is: $`q|kq|k^{}`$. This contradicts the definition of $`\overline{}_a`$.
###### Proposition 0.B.2
Let $`_{T,q}`$ be a box. Let $`s=q\stackrel{u}{}q^{}`$ be an arched execution of $`\overline{}_{T,q}`$ with $`q|k=w`$ and $`q^{}|k=w^{}`$ for all $`kK`$. Then there is a word $`v\mathrm{\Sigma }^{}`$ such that $`v\rho (u)`$ and $`w\stackrel{v}{}w^{}`$ is a path of $`_{T,q}`$.
Proof. We proceed by induction on the size of $`T`$. The case where $`T=\mathrm{}`$ is trivial because $`_{\mathrm{},q}`$ consists of a single state $`q`$. Suppose that the property holds for all subsets $`T^{}T`$ and all states $`q^{}Q`$. Assume first that $`T`$ is a connected set of actions. By Proposition 0.B.1, we know that $`𝕃^{\mathrm{}}(\gamma (s|k))=𝕃^{\mathrm{}}(\gamma (s|k^{}))`$ for all processes $`k,k^{}\mathrm{Loc}(T_1)`$. We claim first that we can find an other execution $`s^{}=q\stackrel{u^{}}{}q^{}`$ such that for *all* processes $`k,k^{}K`$, $`𝕃^{\mathrm{}}(\gamma (s^{}|k))=𝕃^{\mathrm{}}(\gamma (s^{}|k^{}))`$ and moreover $`\rho (u)=\rho (u^{})`$. Let $`𝕃^{\mathrm{}}(\gamma (s^{}|k))=𝒯_1\mathrm{}𝒯_n`$ be the sequence of triangles visited by $`\gamma (s^{}|k)`$. We can split the execution $`s^{}`$ into several smaller arched executions $`s_1,\mathrm{},s_n`$ such that each execution $`s_i`$ is located within triangle $`𝒯_i`$. Similarly each execution $`s_i`$ can be split into several smaller arched executions $`s_1^{},\mathrm{},s_m^{}`$ such that each execution $`s_j^{}`$ is located within a box $`_j`$ inserted in $`𝒯_i`$. Then we can conclude by applying the inductive hypothesis on each smaller box.
Assume finally that $`T`$ is an unconnected set of actions. By Proposition 0.B.1, we know that for all processes $`k,k^{}K`$, $`𝕃^{\mathrm{}}(\gamma (s|k))=𝕃^{\mathrm{}}(\gamma (s|k^{}))`$. Then we can conclude by applying the inductive hypothesis on the smaller boxes visited by $`s`$.
Lemma 3.3 follows now immediately: We have $`\rho (L(\overline{}_{T,q}))[L(_{T,q})]`$. |
warning/0506/math0506053.html | ar5iv | text | # Comparison of semi-simplifications of Galois representations
## Introduction
The aim of this paper is to provide a criterion to ensure that the semi-simplifications of $`p`$-adic finite dimensional Galois representations are isomorphic. Such an isomorphism implies that the Artin L-functions of these representations are the same. This can for example be used to compare the Artin L-functions obtained from automorphic representations with those issued from algebraic geometry. Another possible use is shown in Section 4 which proves that the two representations considered in , one of which is a subrepresentation of the cohomology of a variety while the other is a conjectural automorphic representation, are isomorphic. The main result of this paper, Theorem 3, provides an effective criterion to check whether two semi-simplifications are isomorphic and explains which Frobenius elements suffice to compare the two. In \[3, section 4\], Ron Livné explained and generalized (by lowering the required number of comparisons) a result of Jean-Pierre Serre giving a sufficient condition for semi-simplifications of $`p`$-adic Galois representations to be isomorphic. We intend to generalize here the original result of Jean-Pierre Serre. Even though our result is valid for all dimensions and cannot use the fact that a group of exponent $`2`$ is abelian, our result is similar in complexity to the one of Livné<sup>(1)</sup><sup>(1)</sup>(1)Livné has a better result because he shows that he does not need to compare the representations for all Frobenius elements but only for a so called “non cubic” family.
This paper can be considered as an application of the method explained in \[4, p27–29\].
I would like to thank Bert van Geemen who brought this subject to my attention, gave me some hints and helped remove some errors. This result has been made possible thanks to the help of Thomas Weigel who helped me understand the complexity of the pro-$`p`$ groups and suggested the use of the “powerful pro-$`p`$ groups” which proved to be the right way to tame that complexity. I also would like to thank Karim Belabas who has been of great help for the computational part of the result. I would finally like to thank the referee for his extremely careful reading and his very precise and insightful remarks.
## 1 The result
### 1.1 Setup
This section sets up the framework of this work. Let us fix an integer $`n2`$, a prime $`p`$ and define $`m`$ as the minimum integer such that $`p^mn`$. We fix a global field $`K`$ and let $`\overline{K}`$ be a maximal separable algebraic extension of $`K`$. All extensions of $`K`$ considered in this paper are sub-extensions of $`\overline{K}`$. For any subfield $`L`$ of $`\overline{K}`$, we denote $`\mathrm{\Gamma }_L=\text{Gal}\left(\overline{K}/L\right)`$.
We denote $`(n,A)`$ the algebra of matrices of size $`n\times n`$ with coefficients in a ring $`A`$.
###### Definition 1
Let $`E`$ be a finite extension of $`\text{Q}_p`$ for some prime $`p`$. Let $`M_1`$ and $`M_2`$ be two matrices in $`(n,E)`$. Let $`F`$ be a finite extension of $`E`$ containing the eigenvalues of $`M_1`$ and $`M_2`$. Denote $`𝒪_F`$ the integer ring of $`F`$, $`𝔭_F`$ its maximal ideal and $`\varpi _F`$ a uniformiser. The two matrices $`M_1`$ and $`M_2`$ are said to have congruent eigenvalues if there exist $`\lambda 𝒪_F^\times `$ and $`v\text{Z}`$ such that the characteristic polynomials of $`\varpi _F^vM_1`$ and $`\varpi _F^vM_2`$ are in $`𝒪_F[X]`$ and are congruent to $`(X\lambda )^n`$ modulo $`𝔭_F`$.
###### Remark 2
* The absolute Galois group of a global field is compact, so the eigenvalues of the matrices in the image of a Galois representation necessarily have valuation $`v=0`$.
* The condition on the matrices is rather strong: it implies that the $`2n`$ eigenvalues are congruent to a single one. Therefore the condition is strong even for each matrix separately.
### 1.2 Construction
For any finite set of places $`S`$ of $`K`$, we want to construct an extension $`K_S=K_{S,n}`$ of $`K`$ such that the Galois group of $`K_S/K`$ is sufficient to compare the semi-simplifications of representations of $`\mathrm{\Gamma }_K`$ with values in $`\text{GL}(n,E)`$, unramified outside $`S`$, and with all eigenvalues reducing to a single one in the residual field of $`E`$.
Take $`K_0=K`$. Define $`K_i`$ by induction by taking $`K_{i+1}`$ to be the maximal abelian extension of $`K_i`$ unramified outside $`S`$ and such that $`\text{Gal}(K_{i+1}/K_i)`$ is a direct product of copies of $`\text{F}_p`$. Notice that $`K_i`$ is a Galois extension of $`K`$ at each step $`i`$. Let $`ϵ=0`$ if $`p2`$ and $`ϵ=1`$ if $`p=2`$. Let $`r=N^{2(1+ϵ)}\frac{N(N1)}{2}`$ with $`N=n[E:\text{Q}_p]`$. Let $`\lambda `$ be the minimum integer such that $`2^\lambda r`$. Finally take $`K_S=K_{\lambda +ϵ+m}`$.
### 1.3 Main result
###### Theorem 3
We fix an integer $`n2`$, a prime $`p`$ and define $`m`$ to be the minimum integer such that $`p^mn`$. Let $`K`$ be a global field, $`S`$ a finite set of places of $`K`$ and $`E`$ a finite extension of $`\text{Q}_p`$. We assume that if $`k`$ is the residual field of $`E`$, then $`n`$ and $`|k^\times |`$ are relatively prime. Let $`K_S`$ be the field constructed as in Section 1.2. Fix a set $`T`$ of places of $`K`$, disjoint from $`S`$, such that each maximal cyclic subgroup of $`\text{Gal}(K_S/K)`$ has a generator of the form $`\text{Frob}(𝔱/t)`$ for some $`tT`$ and some prime $`𝔱`$ above $`t`$ in $`K_S`$. Assume now that
$$\rho _1,\rho _2:\mathrm{\Gamma }_K\text{GL}(n,E)$$
are continuous representations unramified outside $`S`$ and satisfy the following conditions:
1. $`\sigma \mathrm{\Gamma }_K`$, $`\rho _1(\sigma )`$ and $`\rho _2(\sigma )`$ have congruent eigenvalues (see Definition 1).
2. $`tT`$, $`\rho _1(\text{Frob}t)`$ and $`\rho _2(\text{Frob}t)`$ have equal characteristic polynomials (where $`\text{Frob}t`$ is any Frobenius element above $`t`$).
Then $`\rho _1`$ and $`\rho _2`$ have isomorphic semi-simplifications.
###### Remark 4
* In this theorem, the condition $`(n,|k^\times |)=1`$ is needed just to ensure that, up to a twist by a character, the residual representations are $`p`$-groups.
* The characteristic polynomial of a matrix $`M`$ of size $`n`$ has coefficients which are symmetric functions of degree at most $`n`$ of the eigenvalues of $`M`$. Over a field of characteristic $`0`$ the sums of the powers of the variables $`(X_i)_{1in}`$ are a basis of the space of symmetric functions in $`(X_i)`$. It follows that there exists a function $`f`$ independent of $`M`$ such that the characteristic polynomial of $`M`$ is $`f(\text{Tr}M,\text{Tr}M^2,\mathrm{},\text{Tr}M^n)`$. Hence we can modify Condition $`(\text{2})`$ above as follows:
either
$$tT,1kn,\text{Tr}\rho _1((\text{Frob}t)^k)=\text{Tr}\rho _2((\text{Frob}t)^k)$$
or
$$tT,\{\begin{array}{cc}1kn1,\hfill & \text{Tr}\rho _1((\text{Frob}t)^k)=\text{Tr}\rho _2((\text{Frob}t)^k)\hfill \\ & det\rho _1(\text{Frob}t)=det\rho _2(\text{Frob}t).\hfill \end{array}$$
* As for the condition $`(n,|k^\times |)=1`$: observe that, if $`n`$ is even, then $`p`$ has to be $`2`$. Observe also that if $`n`$ is a power of $`p`$, or $`k=\text{F}_2`$, then the condition is verified. Finally observe that we can multiply $`n`$ by $`[E:\text{Q}_p]`$ and asume $`E=\text{Q}_p`$. We can thus always apply the theorem if we choose $`p=2`$ (at the cost of enlarging $`n`$, which makes it less interesting because $`K_S`$ and $`T`$ become larger).
## 2 Pro-$`p`$-groups
The main result of this section is Proposition 9 which establishes our result for a pro-$`p`$ group.
### 2.1 The result for pro-$`p`$ groups
###### Definition 5
For a $`p`$-group or pro-$`p`$ group $`G`$, we denote by $`G^\mathrm{\#}`$ the closure of the intersection of the kernels of all group morphisms from $`G`$ to finite groups such that all their elements have order dividing $`p^m`$.
###### Remark 6
1. $`G^\mathrm{\#}`$ is also called $`\text{}_m(G)`$, at least when $`G`$ is finite.
2. $`G^\mathrm{\#}`$ is normal in $`G`$.
3. Observe that $`G^\mathrm{\#}`$ is also the closure of the subgroup generated by $`p^m`$-th powers.
4. In case $`n=p=2`$, the subgroup $`G^\mathrm{\#}`$ is just the Frattini subgroup $`G^{}`$.
5. If $`\rho :GH`$ is a continuous group morphism, $`\rho (G^\mathrm{\#})=\rho (G)^\mathrm{\#}H^\mathrm{\#}`$ with equality if $`\rho `$ is surjective.
6. If $`G_1`$ and $`G_2`$ are groups, then $`(G_1\times G_2)^\mathrm{\#}=G_1^\mathrm{\#}\times G_2^\mathrm{\#}`$.
The following lemma will be useful later on:
###### Lemma 7
Let $`G`$ be a $`p`$-group such that any element of $`G/G^\mathrm{\#}`$ has a representative in $`G`$ of order dividing $`p^m`$. Then $`G^\mathrm{\#}=\{1\}`$.
Proof: Suppose that $`G^\mathrm{\#}\{1\}`$. Observe first that, according to \[5, Theorem 1.12, p90\], we can find a normal subgroup $`N`$ of $`G`$ which is a subgroup of index $`p`$ of $`G^\mathrm{\#}`$. Then $`\left(G/N\right)^\mathrm{\#}G^\mathrm{\#}/N\text{F}_p`$, so that we can as well assume that $`G^\mathrm{\#}=\text{F}_p`$. We have an exact sequence
$$0\text{F}_pGG/G^\mathrm{\#}1$$
such that each element of $`G/G^\mathrm{\#}`$ has a representative in $`G`$ of order dividing $`p^m`$.
Denote $`H=G/G^\mathrm{\#}`$. Then $`H`$ is a $`p`$-group and $`\text{Aut}\text{F}_p`$ has $`p1`$ elements, so that the action of $`H`$ on $`\text{F}_p`$ is trivial. This means that the extension
$$0\text{F}_pG\stackrel{\pi }{}H1$$
is central, i.e. that $`G^\mathrm{\#}Z(G)`$. Thus every element $`g`$ of $`G`$ has order dividing $`p^m`$ (all elements of $`gG^\mathrm{\#}`$ have the order of $`g`$, except if $`gG^\mathrm{\#}`$, in which case all elements have order either $`p`$ or $`1`$). We deduct that the identity is a morphism from $`G`$ to a group having elements of order dividing $`p^m`$. This means that $`G^\mathrm{\#}=\{1\}`$, which is impossible. $`\mathrm{}`$
###### Remark 8
This lemma is a generalization of \[3, Lemma 4.5, p257\]. The definition of $`G^\mathrm{\#}`$ accounts for Remark 4.6.a. below the proof of the lemma in *loc. cit.*.
###### Proposition 9
Let $`G`$ be a pro-$`p`$ group which is topologically finitely generated and let $`E`$ be a finite extension of $`\text{Q}_p`$. Recall that the integer $`m`$ used to define $`G^\mathrm{\#}`$ is the minimum integer such that $`p^mn`$. Assume
$$\rho _1,\rho _2:G\text{GL}(n,E)$$
are continuous representations and $`\mathrm{\Sigma }G`$ is a subset satisfying:
1. the image of $`\mathrm{\Sigma }^{}=\{\sigma ^k/\sigma \mathrm{\Sigma },k\text{N}\}`$ in $`G/G^\mathrm{\#}`$ is equal to $`G/G^\mathrm{\#}`$;
2. $`\sigma \mathrm{\Sigma }`$, $`\rho _1(\sigma )`$ and $`\rho _2(\sigma )`$ have the same characteristic polynomial.
Then $`\rho _1`$ and $`\rho _2`$ have isomorphic semi-simplifications.
Proof: Let $`𝒪`$ be the integer ring of $`E`$. Since $`G`$ is compact, it preserves a full lattice in $`E^n`$ when acting via each $`\rho _i`$, for $`i=1`$, 2. Since $`𝒪`$ is a discrete valuation ring, such a lattice is free over $`𝒪`$. Hence we may assume $`\rho _i(G)\text{GL}(n,𝒪)`$ for each $`i=1`$, $`2`$.
Let $`𝔭`$ be the maximal ideal of $`𝒪`$ and set $`k=𝒪/𝔭`$. The reduction modulo $`𝔭`$ of $`\rho _i(G)`$ is a $`p`$-group in $`\text{GL}(n,k)`$. A $`p`$-Sylow subgroup for $`\text{GL}(n,k)`$ is the subgroup of upper triangular unipotent matrices. We can thus suppose, up to a base change in the lattices above, that the reduction of $`\rho _i(G)`$ modulo $`𝔭`$ is included in this subgroup. In particular, for any $`g`$ in $`G`$, $`(\rho _i(g)I_n)^n0(mod𝔭)`$. We also have that $`\rho _i(g)^{p^m}I_nmod𝔭`$ (in fact we can substitute $`p^m`$ by any power of $`p`$ that is at least equal to the nilpotency order of the reduction mod $`𝔭`$ of $`\rho _i(g)I_n`$).
Now let $`M_n=(n,𝒪)`$. We define $`\rho :GM_n\times M_n`$ to be the map $`\rho (g)=(\rho _1(g),\rho _2(g))`$. Set $`M`$ to be the linear $`𝒪`$-span of $`\rho (G)`$ in $`M_n\times M_n`$. Then $`M`$ is an $`𝒪`$-algebra spanned (as an $`𝒪`$-module) by $`\mathrm{\Gamma }=\rho (G)`$. Let $`R=M/𝔭M`$ and for $`gG`$, we will denote the image of $`\rho (g)`$ in $`R`$ by $`\overline{g}`$. Set $`\overline{\mathrm{\Gamma }}=\{\overline{g}/gG\}`$. Then $`R`$ is a $`k`$-algebra with unity $`\overline{1}=(I_n,I_n)mod𝔭M`$ and spanned by $`\overline{\mathrm{\Gamma }}`$ as a $`k`$-vector space.
We would like to prove that $`R`$ is spanned over $`k`$ by $`\overline{\mathrm{\Sigma }^{}}=\{\overline{\sigma ^k}/\sigma \mathrm{\Sigma },k\text{N}\}`$. We claim that, for any $`\sigma \mathrm{\Sigma }`$, we have $`(\overline{\sigma }\overline{1})^n=\overline{0}`$ and $`\overline{\sigma }^{p^m}=\overline{1}`$. Both these equalities generalize $`\overline{\sigma }^2=\overline{1}`$ for $`p=n=2`$. The point is that equalities in $`\text{GL}(n,k)`$ can sometimes be translated to equalities in $`R`$. Let us first observe that the characteristic polynomial of $`\rho _i(\sigma )mod𝔭`$ is $`(X1)^n`$. This polynomial is the reduction modulo $`𝔭`$ of the characteristic polynomial of $`\rho _i(\sigma )`$. Let $`_{r=0}^nc_{r,i}X^r`$ be the characteristic polynomial of $`\rho _i(\sigma )`$. Let $`a_{r,i}=(1)^{nr}\left(\genfrac{}{}{0pt}{}{n}{r}\right)c_{r,i}`$. Then $`(\rho _i(\sigma )I_n)^n=_{r=0}^na_{r,i}\rho _i(\sigma )^r`$ and all $`a_{r,i}𝔭`$. From Hypothesis 2, we know that the characteristic polynomials are equal and thus $`a_{r,1}=a_{r,2}=a_r`$. We can deduct that
$`(\rho (\sigma )(I_n,I_n))^n`$ $`=((\rho _1(\sigma )I_n)^n,(\rho _2(\sigma )I_n)^n)`$
$`=({\displaystyle \underset{r=0}{\overset{n}{}}}a_r\rho _1(\sigma )^r,{\displaystyle \underset{r=0}{\overset{n}{}}}a_r\rho _2(\sigma )^r)`$
$`={\displaystyle \underset{r=0}{\overset{n}{}}}a_r(\rho _1(\sigma )^r,\rho _2(\sigma )^r)`$
$`={\displaystyle \underset{r=0}{\overset{n}{}}}a_r(\rho _1(\sigma ),\rho _2(\sigma ))^r`$
$`={\displaystyle \underset{r=0}{\overset{n}{}}}a_r\rho (\sigma )^r`$
$`𝔭M`$
Thus $`N(\sigma )=\overline{\sigma }\overline{1}`$ is nilpotent of order (at most) $`n`$. This means that for any $`r`$, $`\overline{\sigma }^r=(\overline{1}+N(\sigma ))^r`$ is a polynomial in $`N(\sigma )`$ of degree at most $`n1`$. For $`r=p^mn`$, then $`\left(\genfrac{}{}{0pt}{}{p^m}{r}\right)`$ will be in $`p\text{Z}𝔭`$ for all $`i[1;p^m1]`$. Thus $`\overline{\sigma }^{p^m}=\overline{1}`$ (as above we can substitute $`p^m`$ by any power of $`p`$ that is at least equal to the nilpotency order of each $`\rho _i(\sigma )I_nmod𝔭`$). In addition, since we have only used the fact that $`\rho _1(\sigma )`$ and $`\rho _2(\sigma )`$ have the same characteristic polynomial, this remains true for all powers of all the elements of $`\mathrm{\Sigma }`$:
$$\sigma \mathrm{\Sigma }^{},\{\begin{array}{cc}\left(\overline{\sigma }\overline{1}\right)^n=\overline{0}\hfill & \\ \overline{\sigma }^{p^m}=\overline{1},\hfill & \end{array}$$
which means
$$\overline{\sigma }\overline{\mathrm{\Sigma }^{}},\{\begin{array}{cc}(\overline{\sigma }\overline{1})^n=\overline{0}\hfill & \\ \overline{\sigma }^{p^m}=\overline{1}.\hfill & \end{array}$$
To prove that $`R`$ is $`k`$-spanned by $`\overline{\mathrm{\Sigma }^{}}`$ we first prove that $`\overline{\mathrm{\Gamma }}^\mathrm{\#}=\{\overline{1}\}`$. Observe that, since $`𝒪`$ is a principal domain, $`R`$ is a finite-dimensional $`k`$-vector space of dimension at most $`2n^2`$. Hence $`R`$ and $`\overline{\mathrm{\Gamma }}`$ are finite. We can apply Lemma 7 to show that $`\overline{\mathrm{\Gamma }}^\mathrm{\#}=\{1\}`$: since $`\rho (G)=\mathrm{\Gamma }`$, we have $`\mathrm{\Gamma }^\mathrm{\#}=\rho (G)^\mathrm{\#}=\rho (G^\mathrm{\#})`$, which implies $`\overline{\mathrm{\Gamma }}^\mathrm{\#}=\overline{\mathrm{\Gamma }^\mathrm{\#}}=\overline{\rho (G^\mathrm{\#})}=\overline{G^\mathrm{\#}}`$ and thus any element of $`\overline{\mathrm{\Gamma }}/\overline{\mathrm{\Gamma }}^\mathrm{\#}`$ can be represented by an element of $`\overline{\mathrm{\Sigma }^{}}`$ and these elements have order dividing $`p^m`$. According to Lemma 7, we have $`\overline{\mathrm{\Gamma }}^\mathrm{\#}=\{1\}`$ and thus $`\overline{\mathrm{\Gamma }}\overline{\mathrm{\Gamma }}/\overline{\mathrm{\Gamma }}^\mathrm{\#}\overline{\mathrm{\Gamma }/\mathrm{\Gamma }^\mathrm{\#}}\overline{G/G^\mathrm{\#}}\overline{\mathrm{\Sigma }^{}}`$ (the last inclusion is up to the canonical projection from $`G`$ to $`G/G^\mathrm{\#}`$); since $`\overline{\mathrm{\Sigma }^{}}\overline{\mathrm{\Gamma }}`$ and both are finite, we conclude that $`\overline{\mathrm{\Sigma }^{}}=\overline{\mathrm{\Gamma }}`$.
Using the former argument, we can apply Nakayama’s lemma to see that $`\mathrm{\Sigma }^{}`$ generates $`M`$ as an $`𝒪`$-module. Since the characteristic polynomials of $`\rho _1(\sigma )`$ and $`\rho _2(\sigma )`$ are equal, the traces of $`\rho _1(\sigma ^k)`$ and $`\rho _2(\sigma ^k)`$ are equal for all $`\sigma \mathrm{\Sigma }`$ and all $`k\text{N}`$. Thus the linear form $`\alpha `$ on $`M`$ defined by $`\alpha (a,b)=\text{Tr}a\text{Tr}b`$ is trivial on a generating set of $`M`$ and thus on all of $`M`$. As a consequence, the characteristic polynomials of $`\rho _1(g)`$ and $`\rho _2(g)`$ are equal for all $`gG`$. $`\mathrm{}`$
### 2.2 Structure of pro-$`p`$ groups
A good reference for the following is , and in particular chapter 3.
###### Definition 10
A powerful pro-$`p`$ group is a pro-$`p`$ group $`G`$ such that $`G/G^p`$ (resp. $`G/G^4`$ if $`p=2`$) is abelian, where $`G^p`$ (resp. $`G^4`$) is the subgroup generated by $`p`$-th (resp. fourth) powers of elements of $`G`$.
###### Proposition 11
For each finitely generated pro-$`p`$ group $`G`$ with a powerful open subgroup, there is a number $`r`$ such that any subgroup of $`G`$ has at most $`r`$ generators.
###### Definition 12
The minimal number $`r`$ above is called the rank of the pro-$`p`$ group $`G`$.
For any integer $`r1`$ we define the integer $`\lambda (r)`$ as the minimum $`\mathrm{}`$ such that $`2^{\mathrm{}}r`$.
A proof of the following result is included in the proof of \[6, Theorem 3.10\].
###### Theorem 13
For any pro-$`p`$ group $`G`$ of rank $`r`$, there exists a $`t\lambda (r)+ϵ`$ and a filtration
$$G_tG_{t1}\mathrm{}G_0=G$$
with abelian quotients of exponent $`p`$ such that $`G_t`$ is powerful. Recall that $`ϵ=1`$ if $`p=2`$ and $`ϵ=0`$ otherwise.
## 3 Reinterpretation of $`G/G^\mathrm{\#}`$ in the Galois group
Proof of Theorem 3: We take $`k`$ to be the residual field of $`E`$ and $`q=|k|`$. Since $`(n,q1)=1`$, the map $`xx^n`$ is injective and thus surjective and bijective in $`k`$. Therefore there exists a unique character
$$\overline{\chi }:\mathrm{\Gamma }_Kk^\times $$
satisfying $`\overline{\chi }^n=det\rho _i(mod𝔭)`$ for $`i=1`$, $`2`$. Let $`\chi `$ be the Teichmüller lift of $`\overline{\chi }`$. Then all the eigenvalues of $`\chi ^1(g)\rho _i(g)`$ will be in some finite extension $`F`$ of $`E`$ and they will reduce to the same $`\lambda `$ in some finite extension $`k^{}`$ of $`k`$. The characteristic polynomial of the reduction mod $`𝔭`$ of each $`\chi ^1(g)\rho _i(g)`$ will be of the form $`P_i(X)=(X\lambda )^n`$. We write $`n=p^vm`$ with $`(m,p)=1`$. We then have
$$P(X)=(X^{p^v}\lambda ^{p^v})^m=X^nm\lambda ^{p^v}X^{np^v}+\mathrm{}+\lambda ^n$$
so that, since $`m0`$ in $`k`$, $`\lambda ^{p^v}k`$. This shows that $`\lambda k`$. Since $`\lambda ^n=\overline{det\rho _i(g)\chi ^1(g)}=1`$, we obtain $`\lambda =1`$. Thus the image of $`\mathrm{\Gamma }_K`$ under the map $`\rho (g)=\chi ^1(g)(\rho _1(g),\rho _2(g))`$ is a pro-$`p`$ group $`G\text{GL}(n,E)^2`$. This can easily be seen from \[6, Proposition 1.11, p22\]: change basis so that both reductions mod $`𝔭`$ of $`\rho _i(\mathrm{\Gamma }_K)`$, for $`i=1`$, $`2`$, have image in the subgroup $`U_k`$ of unipotent upper triangular matrices. Let $`U_𝔭`$ be the inverse image of $`U_k`$ in $`\text{GL}(n,𝒪)`$. Then $`U_k`$ is a $`p`$-group and the kernel of the reduction mod $`𝔭`$ is the normal subgroup $`V=I_n+\varpi (n,𝒪)`$ which is a pro-$`p`$ group. Hence $`U_𝔭`$ is a pro-$`p`$ group and $`\rho _1(\mathrm{\Gamma }_K)`$ and $`\rho _2(\mathrm{\Gamma }_K)`$ are closed in $`U_𝔭`$, because they are compact, thus they also are pro-$`p`$ groups.
We want to compute the ranks of $`G_1=\rho _1(\mathrm{\Gamma }_K)`$ and $`G_2=\rho _2(\mathrm{\Gamma }_K)`$. We begin by embedding $`\text{GL}(n,𝒪)`$ in $`\text{GL}(N,\text{Z}_p)`$ by using a basis of $`𝒪`$ over $`\text{Z}_p`$ to identify $`𝒪^n`$ and $`\text{Z}_p^N`$. Let $`M`$ be a matrix of $`\text{GL}(n,𝒪)`$ with characteristic polynomial $`P(X)`$. Its embedding $`M_r`$ in $`\text{GL}(N,\text{Z}_p)`$ has characteristic polynomial equal to $`P^\sigma (X)`$, where $`\sigma `$ runs over the embeddings of $`E`$ in a fixed algebraic closure of $`E`$ and $`P^\sigma `$ is the polynomial obtained from $`P`$ by applying $`\sigma `$ to its coefficients. In particular, if $`M`$ reduces to an unipotent matrix in $`\text{GL}(n,k)`$, its characteristic polynomial is congruent to $`(X1)^n`$ modulo $`𝔭`$ so that the characteristic polynomial of $`M_r`$ is congruent to $`(X1)^N`$ modulo $`p`$. This means that $`M_r`$ reduces to an unipotent matrix in $`\text{GL}(N,\text{F}_p)`$. The group of unipotent matrices of $`\text{GL}(N,\text{F}_p)`$ has rank at most $`\frac{N(N1)}{2}`$. The kernel of the reduction mod $`p`$ in $`\text{GL}(N,\text{Z}_p)`$ is $`V=I_n+p(N,\text{Z}_p)`$. According to \[6, Theorem 5.2, p88\], if $`p`$ is odd then $`V`$ is powerful of rank $`N^2`$ while if $`p=2`$ then the subgroup $`V^{}=I_n+4(N,\text{Z}_2)`$ is powerful of rank $`N^2`$ and $`V/V^{}`$ is a subgroup of $`(\text{Z}/2\text{Z})^{N^2}`$, which means that it is a $`2`$-group of rank at most $`N^2`$. Putting all three terms together, we see that the group of matrices that reduce to the subgroup of unipotent matrices in $`\text{GL}(n,k)`$ has rank at most $`r=N^2(N^2)^ϵ\frac{N(N1)}{2}`$. This means that the ranks of $`G_1`$ and $`G_2`$ are at most $`r`$. We can apply Theorem 13 to $`G_i`$, for $`i=1`$, $`2`$: for some $`t\lambda (r)+ϵ`$, we get a filtration
$$V_i=G_{i,t}G_{i,t1}\mathrm{}G_{i,1}G_{i,0}=G_i$$
with all quotients $`G_{i,s}/G_{i,s+1}`$ abelian of exponent $`p`$ and $`V_i`$ a powerful pro-$`p`$ group. Since $`V_i`$ is powerful, with $`m`$ more filtration steps we get $`V_i^\mathrm{\#}`$. It is clear that since $`V_iG_i`$, we have $`V_i^\mathrm{\#}G_i^\mathrm{\#}`$. Since $`G^\mathrm{\#}`$ is the closure of the subgroup generated by the $`p^m`$-th powers, we see that $`G^\mathrm{\#}G_1^\mathrm{\#}\times G_2^\mathrm{\#}`$. This means that a filtration with at most $`\lambda (r)+ϵ+m`$ steps is sufficient to get a subgroup $`V^\mathrm{\#}`$ of $`G^\mathrm{\#}`$<sup>(2)</sup><sup>(2)</sup>(2)Observe that $`r`$ is not an upper bound for the rank of $`G`$: the rank of $`G`$ is at most $`2r`$ but can be greater than $`r`$..
On the field side, the $`i`$-th step of the filtration corresponds to an extension of $`K_i`$ by an abelian extension of exponent $`p`$, i.e. the compositum of cyclic extensions of order $`p`$. This means that $`\rho (\mathrm{\Gamma }_{K_S})V^\mathrm{\#}G^\mathrm{\#}`$.
Then Proposition 9 gives the result. $`\mathrm{}`$
###### Proposition 14
Let $`K`$, $`n`$, $`p`$, $`E`$, $`𝒪`$, $`𝔭`$, $`k`$, $`q`$ and $`S`$ be as in Theorem 3 and its proof. Let $`\rho _1`$$`\rho _2:\mathrm{\Gamma }_K\text{GL}(n,E)`$ be two representations unramified outside $`S`$. Let $`K^{}`$ be the compositum of all extensions of $`K`$ unramified outside $`S`$ with degree $`d`$ such that:
* $`d|\mathrm{\#}\text{GL}(n,k)`$
* $`(d,p)=1`$
* $`d\frac{q^n1}{q1}`$.
Denote $`\rho _1^{}`$ and $`\rho _2^{}`$ the respective restrictions of $`\rho _1`$ and $`\rho _2`$ to $`\mathrm{\Gamma }_K^{}`$. Then $`\rho _1^{}`$ and $`\rho _2^{}`$ satisfy Condition 1 of Theorem 3.
Proof: Let $`G`$ be a subgroup of $`\text{GL}(n,k)`$. We consider a flag $`V_0=\{0\}V_1\mathrm{}V_{\mathrm{}}=k^n`$ such that each $`V_i`$ is stable under the action of $`G`$ and the action of $`G`$ on each quotient $`\overline{V_i}=V_i/V_{i1}`$ is irreducible. We will denote $`G_i`$ the image of $`G`$ in $`\text{Hom}(\overline{V_i})`$ and $`d_i=\text{dim}_k\overline{V_i}`$. In a basis adapted to the flag $`(V_i)`$, the matrices representing the action of $`G`$ on $`k^n`$ are blockwise upper-triangular and the $`i`$-th diagonal block of an element $`gG`$ is equal to the projection of $`g`$ in $`G_i`$.
Then for any $`i\{1,\mathrm{},\mathrm{}\}`$, $`G_i`$ is a finite group and a subgroup of a general linear group. If $`P`$ is a $`p`$-Sylow subgroup of $`G_i`$, then the elements of $`P`$ are the elements $`gG_i`$ such that, for certain basis $`(e_j)`$ of $`\overline{V_i}`$, $`j\{1,\mathrm{},d_i\}`$, $`g(e_j)=e_j+_{k<j}\lambda _{k,j}e_k`$. In particular a $`p`$-Sylow of $`G_i`$ fixes at least one vector in $`\overline{V_i}`$. Let $`e_i`$ be such a vector, $`\{e_{i,j}\}_{1in_i}`$ its images under the action of $`G_i`$ (with $`e_{i,1}=e_i`$) and $`H_{i,j}=\text{Fix}_{G_i}(e_{i,j})=\{gG_i/g(e_{i,j})=e_{i,j}\}`$. We see that $`H_{i,j}`$ is a conjugate of $`H_{i,1}=\text{Fix}(e_i)`$ and thus contains a $`p`$-Sylow of $`G_i`$, in particular its index in $`G_i`$ is prime to $`p`$. Let $`H_i=_jH_{i,j}`$. Since $`\overline{V_i}`$ is irreducible under the action of $`G_i`$, for any $`v\overline{V_i}`$ there exists $`(\lambda _j)k^{d_i}`$ such that $`v=_j\lambda _je_{i,j}`$. In particular $`gH_i`$, $`g(v)=v`$ which means that $`H_i=\{1\}`$.
Let $`G_{i,j}`$ be the inverse image of each $`H_{i,j}`$ in $`G`$. Since $`G_{i,j}`$ is a subgroup of $`G`$, we have $`[G:G_{i,j}]|\mathrm{\#}G`$. Since $`G_i`$ is a projection of $`G`$, we also have $`[G:G_{i,j}]=[G_i:H_{i,j}]q^{d_i}1`$ and $`([G:G_{i,j}],p)=([G_i:H_{i,j}],p)=1`$. It is clear that $`G_{i,j}Z(\text{GL}(n,k))=\{1\}`$ because all the elements of $`H_{i,j}`$ have eigenvalues equal to $`1`$. Consider $`G^{}=k^\times G`$ then we also have $`[G^{}:G_{i,j}]q^{d_i}1`$ (because $`G_{i,j}`$ is the fixator of $`e_{i,j}`$ also in $`G^{}`$). This means that $`[G^{}:k^\times G_{i,j}]\frac{q^{d_i}1}{q1}`$. This in turns implies that if $`Z_0=Z(\text{GL}(n,k))G`$ and $`G_{i,j}^{}=Z_0G_{i,j}`$, then $`[G^{}:k^\times G_{i,j}]=[G:G_{i,j}^{}]`$ so that $`[G:G_{i,j}^{}]\frac{q^{d_i}1}{q1}`$. The intersection of all the $`G_{i,j}`$ project trivially in each $`G_i`$, which means that its elements have eigenvalues equal to $`1`$; we thus see that any $`g_{i,j}G_{i,j}^{}`$ have all its eigenvalues equal.
To finish the proof, take $`G=(\rho _1\times \rho _2)(\mathrm{\Gamma }_K)`$ acting on $`k^n\times k^n`$. Observe that $`k^n\times \{0\}`$ and $`\{0\}\times k^n`$ are both stabilized by $`G`$ so that all $`d_in`$. The inverse image in $`\mathrm{\Gamma }_K`$ of $`G_{i,j}^{}`$ defines an extension that have the properties listed in the hypothesis of the proposition. This means that if $`K_1`$ is their compositum, then the elements of $`(\rho _1\times \rho _2)(\mathrm{\Gamma }_{K_1})`$ have all their eigenvalues equal thus, since $`K_1K^{}`$, $`\rho _1^{}`$ and $`\rho _2^{}`$ verify Condition 1 of Theorem 3. $`\mathrm{}`$
## 4 Numerical application
### 4.1 Short version
In , the authors give an example of two non self-dual representations of $`\text{Gal}(\overline{\text{Q}}/\text{Q})`$ (one should note that the representation coming from the automorphic side is only conjectural) and show that they have equal trace for all primes from $`3`$ to $`67`$. We can apply our result to their example. In our terms, we have $`n=3`$, $`p=2`$ (so that $`m=2`$), $`K=\text{Q}`$, $`E=\text{Q}_2[i]`$ and $`S=\{2\text{Z},\mathrm{}\}`$. We denote by $`\rho _1`$ and $`\rho _2`$ the representations they compare. There are no degree $`3`$ and $`7`$ extensions of Q that ramify only in $`S`$ so that, according to Proposition 14, Condition 1 of the theorem is verified. We made a script in gp/pari to search for the extensions described in the construction of the field $`\text{Q}_S`$. We found that the final compositum is a degree $`64`$ field, which we denote $`\text{Q}_{(2)}`$. In the paper , it is shown that the characteristic polynomial of the image of a Frobenius element $`\text{Frob}_𝔭`$ depends only on its trace. As a consequence, all the eigenvalues of $`\rho _i(\text{Frob}_𝔭)`$ are determined by $`\text{Tr}\rho _i(\text{Frob}_𝔭)`$. The eigenvalues of $`\rho _i(\text{Frob}_𝔭^k)=\rho _i(\text{Frob}_𝔭)^k`$ are powers of the eigenvalues of $`\rho _i(\text{Frob}_𝔭)`$, so that the characteristic polynomial of all the $`\rho _i(\text{Frob}_𝔭^k)`$ are determined by $`\text{Tr}\rho _i(\text{Frob}_𝔭)`$. This means that we can restrict the comparison to the traces of the images of the elements generating maximal cyclic subgroups. Thanks to gp/pari, we found a list of primes $`p`$ such that any element of the Galois group of $`\text{Q}_{(2)}`$ over Q is (conjugate to) the power of a Frobenius element above $`p`$. This list is $`\{5,7,11,17,23,31\}`$. The prime $`3`$ is not included just because of the method (and the particular polynomial defining $`\text{Q}_{(2)}`$) used. Observe that all of the primes have already been checked in the paper .
Professor Luis Dieulefait, from Universitat de Barcelona, made me observe that on page 400 of the aforementioned article, the authors note that the geometric representation is absolutely irreducible, which means in particular that it is equal to its semi-simplification. The remark applies obviously also to the conjectural automorphic representation.
###### Corollary 15
The representation and the tentative representation compared in are isomorphic.
Professor Dieulefait also observed that what is said about $`P_5`$ in the aforementioned article is also true for $`P_7=X^3(1+4i)X^2+7(1+4i)X7^3`$ (the field generated by one root of $`P_7`$ is of degree $`6`$ over Q, contains only fourth roots of unity and it is immediate to see that no rational multiple of $`i`$ is a root of $`P_7`$). This means that all the members of the family of $`\mathrm{}`$-adic representations are absolutely irreducible, hence semi-simple.
### 4.2 Longer version
The script used above to look for $`K_S`$ computes the sequence of fields $`(K_i)`$. At each step, it computes linearly independent Kummer extensions of $`K_i`$ and takes their compositum. Since the ramification is rather limited, we tried to detect early (i.e. before computing the compositum) whether an extension is not a sub-extension of $`K_{i+1}`$. For that purpose, we used the fact that the residual extensions are cyclic, therefore we could not have residual extensions larger than $`p^m=4`$. At each step the residual extension is easily computed using class-field theory. We determined that the beginning of the field sequence is as follows: $`K_1`$ is of degree $`4`$ over Q, $`K_2`$ of degree $`32`$ and $`K_3`$ of degree $`64`$. Since $`\text{Gal}(K_3/K_2)`$ is of order $`2`$, $`\text{Gal}(K_S/K_2)`$ is a cyclic subgroup of $`\text{Gal}(K_S/K)`$<sup>(3)</sup><sup>(3)</sup>(3)$`\text{Gal}(K_S/K_3)`$ is the Frattini subgroup of $`\text{Gal}(K_S/K_2)`$ and can be of index $`2`$ if and only if $`\text{Gal}(K_S/K_2)`$ is cyclic of order some power of $`2`$, here at most $`4`$.. Instead of looking for quadratic extensions of $`K_3`$, we checked that cyclic extensions of order $`4`$ of $`K_2`$ all had a too large residual degree, which proved that $`K_3=K_4`$ and thus $`K_S=K_3`$<sup>(4)</sup><sup>(4)</sup>(4)As an additional proof, we checked with gap that there is no order $`128`$ group admitting such a chain of Frattini quotients.. One equation for the extension is
$$\begin{array}{c}x^{64}16x^{61}96x^{60}+144x^{59}+640x^{58}+1424x^{57}+1184x^{56}18960x^{55}41760x^{54}+\hfill \\ \hfill 1376x^{53}+197184x^{52}+686112x^{51}+503136x^{50}361488x^{49}32684x^{48}422688x^{47}+\\ \hfill 3328944x^{46}+194144x^{45}9106992x^{44}+12742688x^{43}13880240x^{42}2172064x^{41}+\\ \hfill 42205032x^{40}81439424x^{39}+70223264x^{38}+5170976x^{37}112924176x^{36}+181443744x^{35}\\ \hfill 120283616x^{34}73923872x^{33}+288559592x^{32}363513856x^{31}+215744096x^{30}+79679200x^{29}\\ \hfill 318677792x^{28}+319483168x^{27}79843680x^{26}217273248x^{25}+333944272x^{24}161711328x^{23}\\ \hfill 190908864x^{22}+496539520x^{21}579760224x^{20}+422942592x^{19}146636736x^{18}98472864x^{17}+\\ \hfill 232483000x^{16}266632896x^{15}+254039136x^{14}234357888x^{13}+215933024x^{12}\\ \hfill 190302336x^{11}+152557600x^{10}108211328x^9+67231888x^836439104x^7+17140160x^6\\ \hfill 6942400x^5+2395872x^4691136x^3+159168x^226240x+2308.\end{array}$$
Its Galois group is identified in Gap’s small group library as $`[64,34]`$. Up to conjugacy, this group has $`6`$ maximal cyclic subgroups. We list them below using the following convention: if a cyclic subgroup is $`\{1,g,g^2,\mathrm{},g^k\}`$, we write it as $`(1,p_1,p_2,\mathrm{},p_k)`$ where $`p_i`$ is a prime number such that there is a Frobenius element above $`p_i`$ that is equal to $`g^i`$. The list is:
$$(1,5,137,13);(1,7,257,7);(1,11,73,19);(1,17,337,17);(1,23,257,23);(1,31)$$
The center of the group is a two element subgroup generated by $`\text{Frob}(337)`$.
Since we have $`K_3=K_4`$ and there are no extensions degree $`3`$, $`5`$, $`7`$, $`9`$ or $`15`$ of Q ramifying only in $`S`$, the discussion above applies also to $`n=4`$.
Hence, to test for the isomorphism of semi-simplification of representations of $`\text{Gal}(\overline{\text{Q}}/\text{Q})`$ of dimension $`3`$ or $`4`$ over any finite extension of $`\text{Q}_2`$ having $`\text{F}_2`$ as residual field, ramifying only at $`2`$ and $`\mathrm{}`$, it is sufficient to either test
* the traces at primes $`\{5,7,11,13,17,19,23,31,73,137,257,337\}`$;
or
* the characteristic polynomials at primes $`\{5,7,11,17,23,31\}`$. |
warning/0506/math0506409.html | ar5iv | text | # Multiscale Homogenization of Convex Functionals with Discontinuous Integrand
## 1. Introduction
Multiscale composites are structures constituted by two or more materials which are finely mixed on many different microscopic scales. The fact that a composite often combines the properties of the constituent materials makes these structures particularly interesting in many fields of science. There is a vast literature on the subject; we refer the reader to and references therein.
Determining macroscopic behavior of these strongly heterogeneous structures when the size $`\epsilon `$ of the heterogeneity becomes “small” is the aim of homogenization theory.
In the particular case of a periodic multiscale composite, from a variational point of view, the homogenization problem is to characterize the behavior, for the parameter $`\epsilon `$ tending to zero, of functionals on $`W^{1,p}(\mathrm{\Omega },^m)`$ of the type
$$F_\epsilon (u)=_\mathrm{\Omega }f(x,\frac{x}{\rho _1(\epsilon )},\mathrm{},\frac{x}{\rho _n(\epsilon )},u(x))𝑑x,$$
(1.1)
where $``$ denotes the fractional part of a vector componentwise, $`\mathrm{\Omega }`$ is an open bounded domain in $`^d`$, $`\mathrm{}`$ is the unit cell $`[0,1)^d`$, $`\rho _k`$ are the length scales and $`f=f(x,y^1,\mathrm{},y^n,z)`$ is a non-negative function on $`\mathrm{\Omega }\times \mathrm{}^n\times 𝕄^{m\times d}`$.
The purpose of this paper is to analyze (1.1) under the following assumptions.
###### Assumption 1.
$`f`$ is convex in the argument $`z`$ for all $`x\mathrm{\Omega }`$ and $`y^1,\mathrm{},y^n\mathrm{}`$.
###### Assumption 2.
$`f`$ is $`p`$-coercive and with $`p`$-growth:
$$c_1\left|z\right|^pf(x,y^1,\mathrm{},y^n,z)c_2\left(1+\left|z\right|^p\right)$$
for some $`p(1,+\mathrm{})`$, $`c_1,c_2>0`$ and for all $`(x,y^1,\mathrm{},y^n,z)\mathrm{\Omega }\times \mathrm{}^n\times 𝕄^{m\times d}`$.
###### Assumption 3.
$`f`$ is an *admissible integrand*, i.e., for every $`\delta >0`$ there exist a compact set $`X\mathrm{\Omega }`$ with $`\left|\mathrm{\Omega }\backslash X\right|\delta `$ and a compact set $`Y\mathrm{}`$ with $`\left|\mathrm{}\backslash Y\right|\delta `$, such that $`f|_{X\times Y^n\times 𝕄^{m\times d}}`$ is continuous.
In particular we cover the following two significant cases (see Examples 4.12 and 4.13).
* The case of a single microscale ($`n=1`$): the function $`f:\mathrm{\Omega }\times \mathrm{}\times 𝕄^{m\times d}[0,+\mathrm{})`$ is continuous in $`x`$, measurable in $`y`$ and satisfies Assumptions 1 and 2. Notice that $`f`$ is continuous in $`z`$ uniformly with respect to $`x`$ and hence is continuous in $`(x,z)`$. It is possible to interchange the regularity conditions on $`f`$ requiring the measurability in $`x`$ and the continuity in $`y`$.
* The case of a multiscale mixture of two materials: the function $`f:\mathrm{\Omega }\times \mathrm{}^n\times 𝕄^{m\times d}[0,+\mathrm{})`$ is of the type
$$\begin{array}{cc}\hfill f(x,y^1,\mathrm{},y^n,z)=\underset{k=1}{\overset{n}{}}\chi _{P_k}(y^k)f_1(x,y^1,\mathrm{},y^n,z)& \\ \hfill +\left[1\underset{k=1}{\overset{n}{}}\chi _{P_k}(y^k)\right]& f_2(x,y^1,\mathrm{},y^n,z),\hfill \end{array}$$
(1.2)
where $`\chi _{P_k}`$ ($`k=1,\mathrm{},n`$) is the characteristic function of a measurable subset $`P_k`$ of $`\mathrm{}`$ and the functions $`f_1,f_2:\mathrm{\Omega }\times \mathrm{}^n\times 𝕄^{m\times d}[0,+\mathrm{})`$ are measurable in $`x`$, continuous in $`(y^1,\mathrm{},y^n)`$ and satisfy Assumptions 1 and 2. The regularity conditions on $`f_1`$ and $`f_2`$ can be replaced by the continuity in $`(x,y^1,\mathrm{},y^{n1})`$ and the measurability in the fastest oscillating variable $`y^n`$.
Problems of the type (1.1) have captured the attention of many authors. For instance, the case of a single microscale
$$_\mathrm{\Omega }f(x,\frac{x}{\epsilon },u(x))𝑑x$$
has been studied by Braides (see and also \[9, Chapter 14\]) under Assumption 1 and requiring in addition a $`p`$-growth condition on the integrand $`f`$ and a uniform continuity in $`x`$, precisely
$$\left|f(x,y,z)f(x^{},y,z)\right|\omega (\left|xx^{}\right|)\left[\alpha (y)+f(x,y,z)\right]$$
(1.3)
for all $`x,x^{}^d`$, $`y\mathrm{}`$ and $`z𝕄^{m\times d}`$, where $`\alpha L^1(\mathrm{})`$ and $`\omega `$ is a continuous positive function with $`\omega (0)=0`$. Recently Baía and Fonseca have studied this problem under Assumpion 2 and requiring continuity in $`(y,z)`$ and measurability in $`x`$.
In (see also \[9, Chapter 22\], and ) Braides and Lukkassen study functionals of the form
$$_\mathrm{\Omega }f(\frac{x}{\epsilon },\mathrm{},\frac{x}{\epsilon ^n},u(x))𝑑x.$$
The authors provide an iterated homogenization formula for functions as in (1.2) with an additional request on the functions $`f_1`$ and $`f_2`$ of a uniform continuity, similar to (1.3), with respect to the slower oscillating variables $`y^1,\mathrm{},y^{n1}`$. The same result is obtained by Fonseca and Zappale but with a continuous function $`f`$ satisfying Assumptions 1 and 2.
Since the variable $`x`$ describes the macroscopic heterogeneity of the constituent materials while the variables $`y^1,\mathrm{},y^n`$ describe the microscopic heterogeneity of the composite structure, it is desirable to have the weakest possible regularity on them. In particular, the oscillating variables should be able to describe the discontinuity on the interfaces between different materials. At any rate the only request that $`f`$ is borelian is not enough to obtain a homogenization formula, as it is shown in Examples 5.10 and 5.11 (see also and ).
In order to weaken the continuity assumptions taken in the works cited above, we approach the problem using the multiscale Young measures as in (see also and ). The peculiarity of our work is the introduction of the concept of *admissible integrand* (Definition 4.10). The crucial point is to extend the lower semicontinuity property (3.3) to this kind of integrand: this is achieved in Theorem 4.14.
The paper is organized as follows. In Section 2 we recall concepts and basic facts about Young measures. In Section 3 we introduce the notion of multiscale convergence in the general framework of multiscale Young measures. In Section 4 we discuss the properties of admissible integrands. By Theorems 4.6 and 4.14 we derive, in Section 5, the upper and lower estimates for the $`\mathrm{\Gamma }(L^p)`$-limit of the family $`F_\epsilon `$ (Lemmas 5.7 and 5.5). Finally, in Section 6, we give an iterated homogenization formula.
## 2. Young measures
We gather briefly in this section some of the main results about Young measures, for more details and proofs we refer the reader to and .
We denote with
* $`D`$ a bounded Lebesgue measurable subset of $`^l`$ ($`l1`$), equipped with the Lebesgue
$`\sigma `$-algebra $`(D)`$;
* $`\left|A\right|`$ the Lebesgue measure of a set $`A(D)`$;
* $`S`$ a locally compact, complete and separable metric space, equipped with
the Borel $`\sigma `$-algebra $`(S)`$;
* $`(D,S)`$ the family of measurable functions $`u:DS`$;
* $`C_0(S)`$ the space $`\{\varphi :S\text{continuous}:\delta >0KS\text{compact}:|\varphi (z)|<\delta \text{for}zSK\}`$, endowed with the supremum norm;
* $`(S)`$ the space of Radon measures on $`S`$;
* $`𝒫(S):=\{\mu (S):\mu 0\text{and}\mu (S)=1\}`$ the set of probability measures on $`S`$;
* $`L^1(D,C_0(S))`$ the Banach space of all measurable maps $`xD\stackrel{\mathit{\varphi }}{}\varphi _xC_0(S)`$ such that the quantity $`\varphi _{L^1}:=_D\varphi _x_{C_0(S)}𝑑x`$ is finite;
* $`L_w^{\mathrm{}}(D,(S))`$ the Banach space of all weak\* measurable maps $`xD\stackrel{𝜇}{}\mu _x(S)`$ such that $`\mu _{L_w^{\mathrm{}}}:=\mathrm{ess}sup_{xD}\mu _x_{(S)}`$ is finite;
* $`𝒴(D,S)`$ the family of all weak\* measurable maps $`\mu :D(S)`$ such that $`\mu _x𝒫(S)`$ a.e. $`xD`$.
###### Remark 2.1.
*
* As it is known, the dual of $`C_0(S)`$ may be identified with the set of $`S`$-valued Radon measures through the duality
$$\mu ,\varphi =_S\varphi 𝑑\mu \mu (S)\mathrm{and}\varphi C_0(S).$$
* A map $`\mu :D(S)`$ is said to be *weak\* measurable* if $`x\mu _x,\varphi `$ is measurable for all $`\varphi C_0(S)`$. In particular $`x\mu _x_{(S)}`$ is measurable.
* More precisely, the elements of $`L^1(D,C_0(S))`$, $`L_w^{\mathrm{}}(D,(S))`$ and $`𝒴(D,S)`$ are equivalence classes of maps that agree a.e.; we usually do not distinguish these maps from their equivalence classes.
* $`L_w^{\mathrm{}}(D,(S))`$ can be identified with the dual of $`L^1(D,C_0(S))`$ through the duality
$$\mu ,\varphi =_D\mu _x,\varphi _x𝑑x\mu L_w^{\mathrm{}}(D,(S))\mathrm{and}\varphi L^1(D,C_0(S)).$$
In the following we will refer to the weak\* topology of $`L_w^{\mathrm{}}(D,(S))`$ as the topology induced by this duality pairing.
* Let $`\widehat{𝒴}(D,S):=\{\varrho (D\times S):\varrho 0\text{and}\varrho (A\times S)=\left|A\right|A(D)\}`$. By the Disintegration Theorem , the map which associates to $`\mu 𝒴(D,S)`$ the measure $`\widehat{\mu }\widehat{𝒴}(D,S)`$ defined by
$$\widehat{\mu }(A):=_D\left(_S\chi _A(x,z)𝑑\mu _x(z)\right)𝑑xA(D\times S)$$
induces a bijection between $`𝒴(D,S)`$ and $`\widehat{𝒴}(D,S)`$. Given a function $`f:D\times S`$$`\widehat{\mu }`$-integrable, it turns out that $`f(x,)`$ is $`\mu _x`$-integrable for a.e. $`xD`$, $`x_Sf(x,z)𝑑\mu _x(z)`$ is integrable and
$$_{D\times S}f𝑑\widehat{\mu }=_D\left(_Sf(x,z)𝑑\mu _x(z)\right)𝑑x;$$
this last equality remains true if $`f`$ is $`(D)(S)`$-measurable and non-negative.
The family $`(D,S)`$ can be embedded in $`L_w^{\mathrm{}}(D,(S))`$ associating to every $`u(D,S)`$ the function
$$x\stackrel{\delta _u}{}\delta _{u(x)},$$
where $`\delta _{u(x)}`$ is the Dirac probability measure concentrated at the point $`u(x)`$.
###### Definition 2.2.
A function $`\mu L_w^{\mathrm{}}(D,(S))`$ is called the *Young measure* generated by the sequence $`u_h`$ if $`\delta _{u_h}\mu `$ in the weak\* topology.
###### Remark 2.3.
This notion makes sense: by the identification of $`L_w^{\mathrm{}}(D,(S))L^1(D,C_0(S))^{}`$ and as a direct consequence of the Banach-Alaoglu theorem, every sequence $`u_h`$ in $`(D,S)`$ admits a subsequence generating a Young measure.
The following result is a “light” version of the Fundamental Theorem on Young Measures.
###### Theorem 2.4.
Let $`u_h`$ be a sequence in $`(D,S)`$ generating a Young measure $`\mu `$ and for which the “tightness condition” is satisfied, i.e.,
$$\delta >0K_\delta S\text{compact}:\underset{h^+}{sup}\left|\{x:u_h(x)K_\delta \}\right|\delta .$$
(2.1)
The following properties hold:
1. $`\mu 𝒴(D,S)`$;
2. if $`f:D\times S[0,+\mathrm{})`$ is a Caratheodory integrand, then
$$\underset{h+\mathrm{}}{lim\; inf}_\mathrm{\Omega }f(x,u_h(x))𝑑x_\mathrm{\Omega }\overline{f}(x)𝑑x$$
where
$$\overline{f}(x):=_Sf(x,z)𝑑\mu _x(z);$$
3. if $`f:D\times S`$ is a Caratheodory integrand and $`f(,u_h())`$ is equi-integrable, then $`f`$ is $`\widehat{\mu }`$-integrable and $`f(,u_h())\overline{f}`$ weakly in $`L^1(\mathrm{\Omega })`$.
We remember that a $`(D)(S)`$-measurable function $`f`$ is a Caratheodory integrand if $`f(x,)`$ is continuous for all $`xD`$.
## 3. Multiscale Young measures
We introduce now the notion of multiscale convergence, an extension of the two-scale convergence carried out by Allaire () in joint work with Briane. We present it in the general framework of multiscale Young measures, following essentially the ideas exposed in , and .
We start presenting an example that does not only show a fine and explicit case of Young measure, but it is a fundamental mainstay in this section. Before we add some new notations:
* $`\mathrm{\Omega }`$ is a bounded open subset of $`^d`$ ($`d1`$), equipped with the Lebesgue $`\sigma `$-algebra $`(\mathrm{\Omega })`$;
* $`\mathrm{}`$ is the unit cell $`[0,1)^d`$, equipped with the Lebesgue $`\sigma `$-algebra $`(\mathrm{})`$;
* $`n`$ is the number of scales, a positive integer;
* $`\rho _1,\mathrm{},\rho _n`$ are positive functions of a parameter $`\epsilon >0`$ which converge to $`0`$ as $`\epsilon `$ does, for which the following *separation of scales hypothesis* is supposed to hold:
$$\underset{\epsilon 0^+}{lim}\frac{\rho _{k+1}(\epsilon )}{\rho _k(\epsilon )}=0k\{1,\mathrm{},n1\};$$
* $`p(1,+\mathrm{})`$ and $`q[1,+\mathrm{}]`$ (unless otherwise stated), moreover $`q^{}`$ is the Hölderian conjugate exponent of $`q`$;
* $`C_c^j(\mathrm{\Omega })`$ stands for the space of $`j`$-differentiable functions in $`\mathrm{\Omega }`$ with compact support;
* $`C_{per}^j(\mathrm{}^k)`$ is the space of the functions $`u=u(y^1,\mathrm{},y^k)`$ in $`C^j((^d)^k)`$ separately $`\mathrm{}`$-periodic in $`y^1,\mathrm{},y^k`$;
* $`W_{per}^{1,q}(\mathrm{}^k)`$ denotes the space of the functions $`u=u(y^1,\mathrm{},y^k)`$ in $`W_{loc}^{1,q}((^d)^k)`$ separately $`\mathrm{}`$-periodic in $`y^1,\mathrm{},y^k`$.
We fix a sequence $`\epsilon _h0^+`$ of values of the parameter $`\epsilon `$.
###### Example 3.1.
We denote by $`T`$ the set $`\mathrm{}`$ equipped with the topological and differential structure of the $`d`$-dimensional torus and with the Borel $`\sigma `$-algebra $`(T)`$; any function on $`T`$ can be identified with its periodic extension to $`^d`$, in particular
$$C(T)=C_0(T)C_{per}(\mathrm{}).$$
We consider the sequence $`v_h:\mathrm{\Omega }\mathrm{}^n`$ defined by
$$v_h(x):=(\frac{x}{\rho _1(\epsilon _h)},\mathrm{},\frac{x}{\rho _n(\epsilon _h)}).$$
(3.1)
Here $``$ denotes the fractional part of a vector componentwise. For our example, we need an auxiliary ingredient concerning weak convergence. It is a particular case of \[13, Proposition 3.3\].
###### Theorem 3.2.
Riemann-Lebesgue lemma: given $`\varphi C_{per}(\mathrm{}^n)`$, define $`\varphi _h(x):=\varphi (v_h(x))`$. Then $`\varphi _h_\mathrm{}^n\varphi (y^1,\mathrm{},y^n)𝑑y^1\mathrm{}𝑑y^n`$ weakly\* in $`L^{\mathrm{}}(\mathrm{\Omega })`$.
As consequence of Riemann-Lebesgue lemma, for all $`\phi L^1(\mathrm{\Omega })`$ and $`\varphi C_{per}(\mathrm{}^n)`$
$$_\mathrm{\Omega }\phi (x)\varphi \left(v_h(x)\right)𝑑x_{\mathrm{\Omega }\times \mathrm{}^n}\phi (x)\varphi (y^1,\mathrm{},y^n)𝑑x𝑑y^1\mathrm{}𝑑y^n.$$
The map $`\phi \varphi `$ that takes every $`x\mathrm{\Omega }`$ into $`\phi (x)\varphi ()C_{per}(\mathrm{}^n)`$ belongs to $`L^1(\mathrm{\Omega },C_{per}(\mathrm{}^n))`$. Since the space $`L^1(\mathrm{\Omega })C_{per}(\mathrm{}^n)`$, defined as the linear closure of $`\{\phi \varphi :\phi L^1(\mathrm{\Omega })\text{and}\varphi C_{per}(\mathrm{}^n)\}`$, is dense in $`L^1(\mathrm{\Omega },C_{per}(\mathrm{}^n))`$, we conclude that $`v_h`$ generates the Young measure $`\mu 𝒴(\mathrm{\Omega },T^n)`$ with
$$\mu _x=_\mathrm{}\mathrm{}^n\text{for }\text{a.e.}x\mathrm{\Omega },$$
where $`_\mathrm{}\mathrm{}^n`$ is the restriction to $`\mathrm{}^n`$ of the Lebesgue measure on $`\left(^d\right)^n`$.
###### Definition 3.3.
Let $`u_h`$ be a sequence in $`L^1(\mathrm{\Omega })`$. The sequence $`u_h`$ is said to be multiscale convergent to a function $`u=u(x,y^1,\mathrm{},y^n)L^1(\mathrm{\Omega }\times \mathrm{}^n)`$ if
$$\begin{array}{cc}\hfill \underset{h+\mathrm{}}{lim}_\mathrm{\Omega }\phi (x)& \varphi (\frac{x}{\rho _1(\epsilon _h)},\mathrm{},\frac{x}{\rho _n(\epsilon _h)})u_h(x)dx\hfill \\ & =_{\mathrm{\Omega }\times \mathrm{}^n}\phi (x)\varphi (y^1,\mathrm{},y^n)u(x,y^1,\mathrm{},y^n)𝑑x𝑑y^1\mathrm{}𝑑y^n\hfill \end{array}$$
for any $`\phi C_c^{\mathrm{}}\left(\mathrm{\Omega }\right)`$ and any $`\varphi C_{per}^{\mathrm{}}(\mathrm{}^n)`$. We simply write $`u_hu`$. A sequence in $`L^1(\mathrm{\Omega },^m)`$ is called multiscale convergent if it is so componentwise.
###### Proposition 3.4.
Let $`u_h`$ be an equi-integrable sequence in $`L^1(\mathrm{\Omega })`$ multiscale convergent to a function $`uL^1(\mathrm{\Omega }\times \mathrm{}^n)`$. Then $`u_h`$ converges weakly to $`u_{\mathrm{}}`$ in $`L^1(\mathrm{\Omega })`$, where
$$u_{\mathrm{}}(x):=_\mathrm{}^nu(x,y^1,\mathrm{},y^n)𝑑y^1\mathrm{}𝑑y^n.$$
###### Proof.
An equi-integrable sequence is sequentially weakly compact in $`L^1`$, therefore it is sufficient to prove that $`u_hu_{\mathrm{}}`$ in distribution. But this is a direct consequence of the definition, taking $`\varphi 1`$ . ∎
Let $`u_h`$ be a bounded sequence in $`L^1(\mathrm{\Omega },^m)`$; we consider the sequence $`w_h:\mathrm{\Omega }\mathrm{}^n\times ^m`$ defined by
$$w_h(x):=(\frac{x}{\rho _1(\epsilon _h)},\mathrm{},\frac{x}{\rho _n(\epsilon _h)},u_h(x)).$$
(3.2)
Suppose that $`w_h`$ generates a Young measure $`\mu `$ (at any rate this is true, up to a subsequence). Thanks to the boundness hypothesis, it can be easily proved that $`w_h`$ satisfies tightness condition (2.1), so $`\mu 𝒴(\mathrm{\Omega },T^n\times ^m)`$. Roughly speaking, by Remark 2.1(v), it is possible to piece together $`\mu `$ in a measure $`\widehat{\mu }\widehat{𝒴}(\mathrm{\Omega },T^n\times ^m)`$. Thanks to Example 3.1, actually $`\widehat{\mu }\widehat{𝒴}(\mathrm{\Omega }\times \mathrm{}^n,^m)`$ and so, by Remark 2.1(v) again, it is possible to dismantle this measure in a new function $`\nu 𝒴(\mathrm{\Omega }\times \mathrm{}^n,^m)`$, called the *multiscale Young measure* generated by $`u_h`$. In particular we have:
###### Theorem 3.5.
Let $`\mu 𝒴(\mathrm{\Omega },T^n\times ^m)`$ be the Young measure generated by $`w_h`$ and let $`\nu 𝒴(\mathrm{\Omega }\times \mathrm{}^n,^m)`$ be the multiscale Young measure generated by $`u_h`$. Then
$$\begin{array}{cc}\hfill _\mathrm{\Omega }\left(_{\mathrm{}^n\times ^m}f(x,y^1,\mathrm{},y^n,z)𝑑\mu _x(y^1,\mathrm{},y^n,z)\right)& dx\hfill \\ & =_{\mathrm{\Omega }\times \mathrm{}^n}\left(_^mf(x,y^1,\mathrm{},y^n,z)𝑑\nu _{(x,y^1,\mathrm{},y^n)}(z)\right)𝑑x𝑑y^1\mathrm{}𝑑y^n\hfill \end{array}$$
for all $`f:\mathrm{\Omega }\times \mathrm{}^n\times ^m`$ $`\widehat{\mu }`$-integrable or non-negative $`(\mathrm{\Omega })(T^n\times ^m)`$-measurable.
The next statement lights up the link between Young measures and multiscale convergence. Sometimes we will use in the sequel the shorter notation $`y:=(y^1,\mathrm{},y^n)`$ .
###### Theorem 3.6.
Let $`u_h`$ be a bounded sequence in $`L^q(\mathrm{\Omega },^m)`$, $`q[1,+\mathrm{})`$, generating a multiscale Young measure $`\nu `$. The following properties hold:
1. the center of mass $`\overline{\nu }`$, defined by
$$\overline{\nu }(x,y^1,\mathrm{},y^n):=_^mz𝑑\nu _{(x,y^1,\mathrm{},y^n)}(z),$$
is in $`L^q(\mathrm{\Omega }\times \mathrm{}^n,^m)`$;
2. if $`u_h`$ is equi-integrable, then $`u_h\overline{\nu }`$;
3. if $`f:\mathrm{\Omega }\times T^n\times ^m[0,+\mathrm{})`$ is a Caratheodory integrand, i.e., $`(\mathrm{\Omega })(T^n\times ^m)`$-measurable and continuous on $`T^n\times ^m`$, then
$$\underset{h+\mathrm{}}{lim\; inf}_\mathrm{\Omega }f(x,w_h(x))𝑑x_{\mathrm{\Omega }\times \mathrm{}^n}\overline{f}(x,y^1,\mathrm{},y^n)𝑑x𝑑y^1\mathrm{}𝑑y^n,$$
(3.3)
where
$$\overline{f}(x,y^1,\mathrm{},y^n):=_^mf(x,y^1,\mathrm{},y^n,z)𝑑\nu _{(x,y^1,\mathrm{},y^n)}(z);$$
4. if $`f:\mathrm{\Omega }\times T^n\times ^m`$ is a Caratheodory integrand and $`f(,w_h())`$ is equi-integrable, then $`f(x,y,)`$ is $`\nu _{(x,y)}`$-integrable for a.e. $`(x,y)\mathrm{\Omega }\times \mathrm{}^n`$, $`\overline{f}`$ is in $`L^1(\mathrm{\Omega }\times \mathrm{}^n)`$ and $`f(,w_h())\overline{f}`$.
###### Proof.
Assertion (iii) is a straight consequence of Theorems 2.4(ii) and 3.5. The integrability properties in assertion (iv) follow by Theorem 2.4(iii) and Remark 2.1(v), by noting that $`\widehat{\mu }=\widehat{\nu }`$. In order to prove the multiscale convergence, fixed $`\phi C_c^{\mathrm{}}\left(\mathrm{\Omega }\right)`$ and $`\varphi C_{per}^{\mathrm{}}(\mathrm{}^n)`$, we define the function $`g:\mathrm{\Omega }\times \mathrm{}^n\times ^m`$ by
$$g(x,y,z):=\phi (x)\varphi (y)f(x,y,z).$$
The function $`g`$ is a Caratheodory integrand on $`\mathrm{\Omega }\times T^n\times ^m`$ and $`g(,w_h())`$ is equi-integrable, thus, by Theorems 2.4(iii) and 3.5,
$$_\mathrm{\Omega }g(x,w_h(x))_\mathrm{\Omega }\left(_{\mathrm{}^n\times ^m}g(x,y,z)𝑑\mu _x(y,z)\right)𝑑x=_{\mathrm{\Omega }\times \mathrm{}^n}\phi (x)\varphi (y)\overline{f}(x,y)𝑑x𝑑y.$$
Assertion (ii) follows by applying (iv) with $`f(x,y,z)=z_j`$, $`j=1,\mathrm{},m`$ . Finally, by Jensen’s inequality and (iii) with $`f(x,y,z)=\left|z\right|^q`$, we obtain assertion (i):
$$_{\mathrm{\Omega }\times \mathrm{}^n}\left|\overline{\nu }(x,y)\right|^q𝑑x𝑑y_{\mathrm{\Omega }\times \mathrm{}^n}\left(_^m\left|z\right|^q𝑑\nu _{(x,y)}(z)\right)𝑑x𝑑y\underset{h+\mathrm{}}{lim\; inf}_\mathrm{\Omega }\left|u_h(x)\right|^q𝑑x<+\mathrm{}.$$
###### Remark 3.7.
Actually assertion (ii) is a compactness result about the multiscale convergence: for every sequence $`u_h`$ equi-integrable in $`L^1`$ or bounded in $`L^p`$, there exists a subsequence $`u_{h_i}`$ which generates a multiscale Young measure and therefore multiscale convergent. Remember that a bounded sequence in $`L^p`$ is equi-integrable by Hölder’s inequality.
We conclude with a basic result about bounded sequences in $`W^{1,p}(\mathrm{\Omega })`$.
###### Theorem 3.8.
Given a sequence $`u_h`$ weakly convergent to $`u`$ in $`W^{1,p}(\mathrm{\Omega })`$, we have that $`u_hu`$ and
$$u_hu+\underset{k=1}{\overset{n}{}}_{y^k}u_k$$
for $`n`$ suitable functions $`u_k(x,y^1,\mathrm{},y^k)L^p(\mathrm{\Omega }\times \mathrm{}^{k1},W_{per}^{1,p}(\mathrm{}))`$ .
The proof can be found in \[2, Theorem 2.6\] and in \[4, Theorem 1.6\]. In the first reference, the idea is to work on the image of $`W^{1,2}(\mathrm{\Omega })`$ under the gradient mapping, by characterizing it as the space orthogonal to all divergence-free functions. Instead in the second reference it is used its characterization as the space of all rotation-free fields. This last method is simpler and works for general $`p`$, even if only the case $`p=2`$ is examined in the original statement. Another proof can be found in .
## 4. Continuity results
As first result of this section, we show that it is possible to use in the multiscale convergence a more complete system of “test functions”, not merely $`\psi (x,y)=\phi (x)\varphi (y)`$ with $`\phi C_c^{\mathrm{}}\left(\mathrm{\Omega }\right)`$ and $`\varphi C_{per}^{\mathrm{}}(\mathrm{}^n)`$. Following Valadier , we introduce opportune classes of functions.
###### Definition 4.1.
A function $`\psi :\mathrm{\Omega }\times \mathrm{}^n`$ is said to be *admissible* if there exist a family $`\{X_\delta \}_{\delta >0}`$ of compact subsets of $`\mathrm{\Omega }`$ and a family $`\{Y_\delta \}_{\delta >0}`$ of compact subsets of $`\mathrm{}`$ such that $`\left|\mathrm{\Omega }\backslash X_\delta \right|\delta `$, $`\left|\mathrm{}\backslash Y_\delta \right|\delta `$ and $`\psi |_{X_\delta \times Y_\delta ^n}`$ is continuous for every $`\delta >0`$.
###### Remark 4.2.
It is not restrictive to suppose that the families $`\{X_\delta \}_{\delta >0}`$ and $`\{Y_\delta \}_{\delta >0}`$ are decreasing, i.e., $`\delta ^{}\delta `$ implies $`X_\delta X_\delta ^{}`$ and $`Y_\delta Y_\delta ^{}`$. Otherwise, it is sufficient to consider the new families $`\{\stackrel{~}{X}_\delta \}_{\delta >0}`$ and $`\{\stackrel{~}{Y}_\delta \}_{\delta >0}`$, where
$$\stackrel{~}{X}_\delta :=\underset{ii_\delta }{}X_{2^i},\stackrel{~}{Y}_\delta :=\underset{ii_\delta }{}Y_{2^i}$$
and $`i_\delta `$ is the minimum positive integer such that $`2^{1i_\delta }\delta `$.
Admissible functions have good measurability properties, as stated in the following lemma. We omit the easy proof.
###### Lemma 4.3.
If $`\psi :\mathrm{\Omega }\times \mathrm{}^n`$ is an admissible function, then there exist a set $`X\mathrm{\Omega }`$ with $`\left|\mathrm{\Omega }\backslash X\right|=0`$ and a set $`Y\mathrm{}`$ with $`\left|\mathrm{}\backslash Y\right|=0`$, such that $`\psi |_{X\times Y^n}`$ is borelian. In particular, for every fixed $`\epsilon `$, the function $`x\psi (x,\frac{x}{\rho _1(\epsilon )},\mathrm{},\frac{x}{\rho _n(\epsilon )})`$ is measurable.
###### Definition 4.4.
An admissible function $`\psi `$ is said to be $`q`$-*admissible*, and we write $`\psi 𝒜dm^q`$, if there exists a positive function $`\alpha L^q(\mathrm{\Omega })`$ such that
$$\left|\psi (x,y)\right|\alpha (x)(x,y)\mathrm{\Omega }\times \mathrm{}^n.$$
The next theorem proves that it is possible to use $`𝒜dm^q`$ as system of test functions. The proof is very close to \[26, Proposition 5\]. Before we state the following lemma, that can be derived by \[13, Lemma 3.1\] (see also \[2, Remark 2.13\]). We use the same definition of $`v_h`$ given in (3.1).
###### Lemma 4.5.
Let $`A_k`$ be a measurable subset of $`\mathrm{}`$ for $`k=1,\mathrm{},n`$ and let $`A:=_{k=1}^nA_k`$. Denoted with $`\chi _A`$ the characteristic function of $`A`$, the sequence $`\chi _A(v_h())`$ converges weakly\* to $`\left|A\right|`$ in $`L^{\mathrm{}}(\mathrm{\Omega })`$.
###### Theorem 4.6.
Let $`u_h`$ be a bounded sequence in $`L^q(\mathrm{\Omega })`$, $`q(1,+\mathrm{}]`$, generating a multiscale Young measure $`\nu `$ and let $`\psi 𝒜dm^q^{}`$. Then
$$\underset{h+\mathrm{}}{lim}_\mathrm{\Omega }\psi (x,v_h(x))u_h(x)𝑑x=_{\mathrm{\Omega }\times \mathrm{}^n}\psi (x,y)\overline{\nu }(x,y)𝑑x𝑑y.$$
In particular, taking $`u_h1`$, we obtain
$$\underset{h+\mathrm{}}{lim}_\mathrm{\Omega }\psi (x,v_h(x))𝑑x=_{\mathrm{\Omega }\times \mathrm{}^n}\psi (x,y)𝑑x𝑑y.$$
(4.1)
###### Proof.
Let $`\delta >0`$; by Lusin theorem applied to $`\alpha `$ and by definition of admissible function, there exist two compact sets $`X\mathrm{\Omega }`$ and $`Y\mathrm{}`$ such that $`\left|\mathrm{\Omega }\backslash X\right|\delta `$, $`\left|\mathrm{}\backslash Y\right|\delta `$ and $`\psi |_{X\times Y^n}`$, $`\alpha |_X`$ are continuous. Let $`M:=\mathrm{max}_X\alpha `$; by Tietze-Urysohn’s theorem, $`\psi |_{X\times Y^n}`$ can be extended to a continuous function $`\psi _0`$ on $`\mathrm{\Omega }\times T^n`$ with $`\left|\psi _0(x,y)\right|M`$ for every $`(x,y)\mathrm{\Omega }\times \mathrm{}^n`$. We define on $`\mathrm{\Omega }\times \mathrm{}^n\times `$ the functions
$$f(x,y,z):=\psi (x,y)z\text{and}f_0(x,y,z):=\psi _0(x,y)z.$$
With the same definition of $`w_h`$ given in (3.2), the sequence $`f_0(,w_h())`$ is equi-integrable because $`\left|f_0(x,w_h(x))\right|M\left|u_h(x)\right|`$. By Theorem 3.6(iv), $`f_0(,w_h())\overline{f_0}`$ and therefore, by Proposition 3.4, $`f_0(,w_h())_\mathrm{}^n\psi _0(,y)\overline{\nu }(,y)𝑑y`$ weakly in $`L^1(\mathrm{\Omega })`$ . This is sufficient to assert that
$$\underset{h+\mathrm{}}{lim}_Xf_0(x,w_h(x))𝑑x=_{X\times \mathrm{}^n}\psi _0(x,y)\overline{\nu }(x,y)𝑑x𝑑y.$$
Now
$$\begin{array}{cc}\hfill |_{\mathrm{\Omega }\times \mathrm{}^n}\psi (x,y)\overline{\nu }(x,y)dxdy& _\mathrm{\Omega }f(x,w_h(x))dx||_{(\mathrm{\Omega }X)\times \mathrm{}^n}\psi (x,y)\overline{\nu }(x,y)dxdy|\hfill \\ \hfill +& \left|_{X\times \mathrm{}^n}\left[\psi (x,y)\psi _0(x,y)\right]\overline{\nu }(x,y)𝑑x𝑑y\right|\hfill \\ \hfill +& \left|_{X\times \mathrm{}^n}\psi _0(x,y)\overline{\nu }(x,y)𝑑x𝑑y_Xf_0(x,w_h(x))𝑑x\right|\hfill \\ \hfill +& \left|_X\left[f_0(x,w_h(x))f(x,w_h(x))\right]𝑑x\right|+\left|_{\mathrm{\Omega }X}f(x,w_h(x))𝑑x\right|\hfill \\ \hfill =& \text{I+II+III+IV+V}.\hfill \end{array}$$
We have to show that I, II, IV and V can be made arbitrarily small. Observe that the function
$$\gamma _\text{I}(x,y):=\alpha (x)_{}\left|z\right|𝑑\nu _{(x,y)}(z)$$
is in $`L^1(\mathrm{\Omega }\times \mathrm{}^n)`$ as consequence of Theorem 3.6(iii), Hölder’s inequality and the $`L^q`$-boundness of the sequence $`u_h`$:
$$\begin{array}{cc}\hfill _{\mathrm{\Omega }\times \mathrm{}^n}\gamma _\text{I}(x,y)𝑑x𝑑y& \underset{h+\mathrm{}}{lim\; inf}_\mathrm{\Omega }\alpha (x)\left|u_h(x)\right|𝑑x\hfill \\ \hfill & \alpha _{L^q^{}(\mathrm{\Omega })}\underset{h}{sup}u_h_{L^q(\mathrm{\Omega })}<+\mathrm{}.\hfill \end{array}$$
The same for $`\gamma _{\text{II}}(x,y):=_\mathrm{\Omega }\left|z\right|𝑑\nu _{(x,y)}(z)`$. By the absolute continuity of the integral and by the estimates
$`\text{I}`$ $`{\displaystyle _{(\mathrm{\Omega }\backslash X)\times \mathrm{}^n}}\gamma _\text{I}(x,y)𝑑x𝑑y`$
and
$`\text{II}=`$ $`{\displaystyle _{X\times (\mathrm{}^n\backslash Y^n)}}\left[\psi (x,y)\psi _0(x,y)\right]\overline{\nu }(x,y)𝑑x𝑑y`$
$``$ $`{\displaystyle _{X\times (\mathrm{}^n\backslash Y^n)}}\gamma _\text{I}(x,y)𝑑x𝑑y+M{\displaystyle _{X\times (\mathrm{}^n\backslash Y^n)}}\gamma _{\text{II}}(x,y)𝑑x𝑑y,`$
we obtain that I and II tend to $`0`$ for $`\delta 0`$. By using again Hölder’s inequality and the $`L^q`$-boundness of $`u_h`$, we get for a suitable positive constant $`c`$
IV $`\left|{\displaystyle _X}\chi _{\mathrm{}^n\backslash Y^n}(v_h(x))\left[f_0(x,w_h(x))f(x,w_h(x))\right]𝑑x\right|`$
$`{\displaystyle _X}\chi _{\mathrm{}^n\backslash Y^n}(v_h(x))\left[\alpha (x)+M\right]\left|u_h(x)\right|𝑑x`$
$`c\left({\displaystyle _X}\chi _{\mathrm{}^n\backslash Y^n}(v_h(x))\left[\alpha (x)+M\right]^q^{}𝑑x\right)^{\frac{1}{q^{}}}`$
and
V $`c\left({\displaystyle _{\mathrm{\Omega }\backslash X}}\left[\alpha (x)\right]^q^{}𝑑x\right)^{\frac{1}{q^{}}}.`$
By Lemma 4.5, it follows that $`\chi _{\mathrm{}^n\backslash Y^n}(v_h())=1\chi _{Y^n}(v_h())`$ converges weakly\* to $`\left|\mathrm{}^n\backslash Y^n\right|`$ and therefore
$$_X\chi _{\mathrm{}^n\backslash Y^n}(v_h(x))\left[\alpha (x)+M\right]^q^{}𝑑x\stackrel{h\mathrm{}}{}\left|\mathrm{}^n\backslash Y^n\right|_X\left[\alpha (x)+M\right]^q^{}𝑑x.$$
Hence we conclude that IV and V tend to $`0`$ for $`h\mathrm{}`$ and $`\delta 0`$ . ∎
###### Remark 4.7.
Let $`\psi =\psi (x,y^1,\mathrm{},y^n)`$ be a real function on $`\mathrm{\Omega }\times \mathrm{}^n`$ either continuous in $`(y^1,\mathrm{},y^n)`$ and measurable in $`x`$ or continuous in $`(x,y^1,\mathrm{},y^{k1},y^{k+1},\mathrm{},y^n)`$ and measurable in $`y^k`$. By Scorza-Dragoni theorem (see ), $`\psi `$ is an admissible function. This is no longer true if one removes the continuity assumption on two variables. More generally, the invocation of (4.1) may be invalid, as shown in the next two examples. The first covers the case $`\psi =\psi (x,y)`$ ($`n=1`$) while the second covers the case $`\psi =\psi (y^1,y^2)`$. We remark that in both examples $`\psi `$ is a Borel function. See also the example in \[1, Proposition 5.8\].
###### Example 4.8.
Define the Borel sets $`A_i:=_{j=0}^{i1}\{(x,y)[0,1)^2:y=ixj\}`$ and $`A:=_{i=1}^{\mathrm{}}A_i`$. Now, in the simple case $`d=n=1`$, $`\rho _1(\epsilon _h)=h^1`$ and $`\mathrm{\Omega }=(0,1)`$, consider the function $`\psi (x,y):=\chi _A(x,y)`$. We have
$$_0^1\psi (x,hx)𝑑x1\text{but}_0^1_0^1\psi (x,y)𝑑x𝑑y=0.$$
###### Example 4.9.
In the case $`d=1`$, $`n=2`$, $`\rho _1(\epsilon _h)=h^1`$, $`\rho _2(\epsilon _h)=h^2`$ and $`\mathrm{\Omega }=(0,1)`$, consider the function $`\psi (y^1,y^2):=\chi _A(y^1,y^2)`$, where $`A`$ is defined as in the former example. We have
$$_0^1\psi (hx,h^2x)𝑑x1\text{but}_0^1_0^1_0^1\psi (y^1,y^2)𝑑x𝑑y^1𝑑y^2=0.$$
Notice that the result of weak\* convergence in $`L^{\mathrm{}}`$ stated in Lemma 4.5 is not applicable to $`A`$.
So far we have considered Caratheodory functions $`f`$ on $`\mathrm{\Omega }\times T^n\times ^m`$. As we explained in the introduction, one would like to have a minimal regularity in $`(x,y^1,\mathrm{},y^n)`$. For this reason, we introduce an opportune class of integrands and extend to this Theorem 3.6(iii).
###### Definition 4.10.
A function $`f:\mathrm{\Omega }\times \mathrm{}^n\times ^m[0,+\mathrm{})`$ is said to be an *admissible integrand* if for every $`\delta >0`$ there exist a compact set $`X\mathrm{\Omega }`$ with $`\left|\mathrm{\Omega }\backslash X\right|\delta `$ and a compact set $`Y\mathrm{}`$ with $`\left|\mathrm{}\backslash Y\right|\delta `$, such that $`f|_{X\times Y^n\times ^m}`$ is continuous.
As in the analogous case for admissible functions (Lemma 4.3), it is easy to verify the following measurability properties of admissible integrands.
###### Lemma 4.11.
If $`f:\mathrm{\Omega }\times \mathrm{}^n\times ^m[0,+\mathrm{})`$ is an admissible integrand, then there exist a set $`X\mathrm{\Omega }`$ with $`\left|\mathrm{\Omega }\backslash X\right|=0`$ and a set $`Y\mathrm{}`$ with $`\left|\mathrm{}\backslash Y\right|=0`$, such that $`f|_{X\times Y^n\times ^m}`$ is borelian. In particular, for every fixed $`\epsilon `$, the function $`(x,z)f(x,\frac{x}{\rho _1(\epsilon )},\mathrm{},\frac{x}{\rho _n(\epsilon )},z)`$ is $`(\mathrm{\Omega })(^m)`$-measurable.
###### Example 4.12.
Let $`f:\mathrm{\Omega }\times \mathrm{}\times ^m[0,+\mathrm{})`$ be a function such that
* $`f(,y,)`$ is continuous for all $`y\mathrm{}`$;
* $`f(x,,z)`$ is measurable for all $`x\mathrm{\Omega }`$ and $`z^m`$.
By Scorza-Dragoni theorem, $`f`$ is an admissible integrand. Clearly it is possible to replace conditions (i) and (ii) with
* $`f(x,,)`$ is continuous for all $`x\mathrm{\Omega }`$;
* $`f(,y,z)`$ is measurable for all $`y\mathrm{}`$ and $`z^m`$.
###### Example 4.13.
Let $`f:\mathrm{\Omega }\times \mathrm{}^n\times ^m[0,+\mathrm{})`$ be a function of the type
$$\begin{array}{cc}\hfill f(x,y^1,\mathrm{},y^n,z)=\underset{k=1}{\overset{n}{}}\chi _{P_k}(y^k)f_1(x,y^1,\mathrm{},y^n,z)& \\ \hfill +\left[1\underset{k=1}{\overset{n}{}}\chi _{P_k}(y^k)\right]& f_2(x,y^1,\mathrm{},y^n,z),\hfill \end{array}$$
where $`P_k`$ ($`k=1,\mathrm{},n`$) is a measurable subset of $`\mathrm{}`$ and $`f_j`$ ($`j=1,2`$) is a non-negative function on $`\mathrm{\Omega }\times \mathrm{}^n\times ^m`$ such that
* $`f_j`$ is continuous in $`(y^1,\mathrm{},y^n,z)`$;
* $`f_j`$ is measurable in $`x`$.
By Scorza-Dragoni theorem for every $`\delta >0`$ there exists a compact set $`X\mathrm{\Omega }`$ such that the functions $`f_1`$ and $`f_2`$ are continuous on $`X\times \mathrm{}^n\times ^m`$. By applying Lusin theorem to each $`\chi _{P_k}`$, we obtain that $`f`$ is an admissible integrand. Obviously, the conditions (i) and (ii) can be replaced by
* $`f_j`$ is continuous in $`(x,y^1,\mathrm{},y^{k1},y^{k+1},\mathrm{},y^n,z)`$;
* $`f_j`$ is measurable in $`y^k`$.
###### Theorem 4.14.
Let $`u_h`$ be a bounded sequence in $`L^q(\mathrm{\Omega },^m)`$, $`q[1,+\mathrm{})`$, generating a multiscale Young measure $`\nu `$ and let $`f:\mathrm{\Omega }\times \mathrm{}^n\times ^m[0,+\mathrm{})`$ be an admissible integrand satisfying the $`q`$-growth condition
$$f(x,y^1,\mathrm{},y^n,z)c\left(1+\left|z\right|^q\right)$$
for some $`c>0`$ and for all $`(x,y^1,\mathrm{},y^n,z)\mathrm{\Omega }\times \mathrm{}^n\times ^m`$ . Then
$$\underset{h+\mathrm{}}{lim\; inf}_\mathrm{\Omega }f(x,w_h(x))𝑑x_{\mathrm{\Omega }\times \mathrm{}^n}\overline{f}(x,y^1,\mathrm{},y^n)𝑑x𝑑y^1\mathrm{}𝑑y^n,$$
(4.2)
where as usual
$$\overline{f}(x,y^1,\mathrm{},y^n):=_^mf(x,y^1,\mathrm{},y^n,z)𝑑\nu _{(x,y^1,\mathrm{},y^n)}(z).$$
###### Proof.
Assume initially that, in addition,
$$f(x,y,z)=0\text{if}\left|z\right|r$$
(4.3)
for a fixed $`r>0`$. By the admissibility condition, for every $`\delta >0`$ there exist a compact set $`X\mathrm{\Omega }`$ and a compact set $`Y\mathrm{}`$ such that $`\left|\mathrm{\Omega }\backslash X\right|\delta `$, $`\left|\mathrm{}\backslash Y\right|\delta `$ and $`f|_{X\times Y^n\times ^m}`$ is continuous. By Tietze-Urysohn’s theorem, $`f|_{X\times Y^n\times ^m}`$ can be extended to a continuous function $`f_0`$ on $`\mathrm{\Omega }\times T^n\times ^m`$ with $`0f_0(x,y,z)M`$ for every $`(x,y,z)\mathrm{\Omega }\times \mathrm{}^n\times ^m`$, where $`M:=\mathrm{max}f`$ on $`X\times Y^n\times ^m`$. Notice that, by the $`q`$-growth condition, $`Mc\left(1+r^q\right)`$. Obviously $`f_0(,w_h())`$ is equi-integrable and so, by Theorem 3.6(iv) and by Proposition 3.4,
$$\underset{h+\mathrm{}}{lim}_\mathrm{\Omega }f_0(x,w_h(x))𝑑x=_{\mathrm{\Omega }\times \mathrm{}^n}\overline{f_0}(x,y)𝑑x𝑑y.$$
For a suitable subsequence $`h_i`$
$$\underset{i+\mathrm{}}{lim}_\mathrm{\Omega }f(x,w_{h_i}(x))𝑑x=\underset{h+\mathrm{}}{lim\; inf}_\mathrm{\Omega }f(x,w_h(x))𝑑x.$$
We can write
$$\begin{array}{cc}\hfill \underset{i+\mathrm{}}{lim}_\mathrm{\Omega }f(x,w_{h_i}(x))𝑑x& _{\mathrm{\Omega }\times \mathrm{}^n}\overline{f}(x,y)𝑑x𝑑y=\underset{i+\mathrm{}}{lim}_\mathrm{\Omega }\left[f(x,w_{h_i}(x))f_0(x,w_{h_i}(x))\right]𝑑x\hfill \\ & +[lim_{i+\mathrm{}}_\mathrm{\Omega }f_0(x,w_{h_i}(x))dx_{\mathrm{\Omega }\times \mathrm{}^n}\overline{f_0}(x,y)dydx]\hfill \\ & +_{\mathrm{\Omega }\times \mathrm{}^n}\left[\overline{f_0}(x,y)\overline{f}(x,y)\right]𝑑x𝑑y=\text{I+II+III}.\hfill \end{array}$$
Let us check that the negative part of I and III can be made arbitrarily small.
Firstly, by Lemma 4.5,
$$\begin{array}{cc}\hfill \text{I}=& \underset{i+\mathrm{}}{lim}_{\mathrm{\Omega }\backslash X}\left[f(x,w_{h_i}(x))f_0(x,w_{h_i}(x))\right]𝑑x\hfill \\ & +\underset{i+\mathrm{}}{lim}_X\chi _{\mathrm{}^n\backslash Y^n}(v_{h_i}(x))\left[f(x,w_{h_i}(x))f_0(x,w_{h_i}(x))\right]𝑑x\hfill \\ \hfill & M\left|\mathrm{\Omega }\backslash X\right|\underset{i+\mathrm{}}{lim}M_X\chi _{\mathrm{}^n\backslash Y^n}(v_{h_i}(x))𝑑x\hfill \\ \hfill & M\left(\left|\mathrm{\Omega }\backslash X\right|+\left|X\right|\left|\mathrm{}^n\backslash Y^n\right|\right)c\left(1+r^q\right)\left(\delta +n\left|\mathrm{\Omega }\right|\delta \right).\hfill \end{array}$$
Now, observe that the function
$$\gamma (x,y):=_^m\left[f(x,y,z)f_0(x,y,z)\right]𝑑\nu _{(x,y)}(z)$$
is in $`L^1(\mathrm{\Omega }\times \mathrm{}^n)`$ as consequence of Theorem 3.6(iii):
$$\begin{array}{cc}\hfill _{\mathrm{\Omega }\times \mathrm{}^n}\left|\gamma (x,y)\right|𝑑x𝑑y& _{\mathrm{\Omega }\times \mathrm{}^n}\left[M+c+c_^m\left|z\right|^q𝑑\nu _{(x,y)}(z)\right]𝑑x𝑑y\hfill \\ \hfill & (M+c)\left|\mathrm{\Omega }\right|+c\underset{h+\mathrm{}}{lim\; inf}_\mathrm{\Omega }\left|u_h(x)\right|^q𝑑x<+\mathrm{}.\hfill \end{array}$$
By the absolute continuity of the integral and by the equality
$$\text{III}=_{(\mathrm{\Omega }\times \mathrm{}^n)(X\times Y^n)}\gamma (x,y)𝑑x𝑑y,$$
we obtain that III tends to $`0`$ for $`\delta 0`$. This concludes the first part of the proof.
In order to remove assumption (4.3) we consider, for $`k^+`$, the functions $`p_kC_0(^m)`$ defined by
$$p_k\left(z\right):=\{\begin{array}{cc}1\hfill & \text{if}\left|z\right|k\hfill \\ 1+k\left|z\right|\hfill & \text{if}k\left|z\right|k+1\hfill \\ 0\hfill & \text{if}\left|z\right|k+1\hfill \end{array}$$
and the functions $`f_k(x,y,z):=p_k(z)f(x,y,z)`$ . By applying the first part of the theorem, we have
$$\underset{h+\mathrm{}}{lim\; inf}_\mathrm{\Omega }f(x,w_h(x))𝑑x\underset{h+\mathrm{}}{lim\; inf}_\mathrm{\Omega }f_k(x,w_h(x))𝑑x_{\mathrm{\Omega }\times \mathrm{}^n}\overline{f_k}(x,y)𝑑x𝑑y.$$
By noting that $`f_k`$ is increasing and that $`f_k(x,y,)f(x,y,)`$ a.e. in $`^m`$ for every fixed $`(x,y)\mathrm{\Omega }\times \mathrm{}^n`$, we deduce from the monotone convergence theorem that $`\overline{f_k}\overline{f}`$ a.e. in $`\mathrm{\Omega }\times \mathrm{}^n`$. The sequence $`\overline{f_k}`$ is increasing so, again from monotone convergence theorem,
$$_{\mathrm{\Omega }\times \mathrm{}^n}\overline{f_k}(x,y)𝑑x𝑑y\stackrel{k\mathrm{}}{}_{\mathrm{\Omega }\times \mathrm{}^n}\overline{f}(x,y)𝑑x𝑑y.$$
###### Remark 4.15.
Lower semicontinuity property (4.2) is not true if $`f`$ is only borelian. For instance, consider the function $`f(x,y,z):=[1\psi (x,y)]\left|z\right|^p`$, where $`\psi `$ is defined as in Example 4.8.
## 5. Gamma-convergence
In the present section we examine the multiperiodic homogenization of nonlinear convex functionals by means of the $`\mathrm{\Gamma }`$-convergence combined with the multiscale Young measures.
Before we recall the definition of $`\mathrm{\Gamma }`$-convergence, referring to and for an exposition of the main properties.
###### Definition 5.1.
Let $`(U,\tau )`$ be a topological space satisfying the first countability axiom and $`F_h`$, $`F`$ functionals from $`U`$ to $`[\mathrm{},+\mathrm{}]`$; we say that $`F`$ is the $`\mathrm{\Gamma }(\tau )`$-limit of the sequence $`F_h`$ or that $`F_h`$ $`\mathrm{\Gamma }(\tau )`$-converges to $`F`$, and write
$$F=\mathrm{\Gamma }(\tau )\text{-}\underset{h+\mathrm{}}{lim}F_h,$$
if for every $`uU`$ the following conditions are satisfied:
$$F(u)inf\{\underset{h+\mathrm{}}{lim\; inf}F_h(u_h):u_h\stackrel{𝜏}{}u\}$$
(5.1)
and
$$F(u)inf\{\underset{h+\mathrm{}}{lim\; sup}F_h(u_h):u_h\stackrel{𝜏}{}u\}.$$
(5.2)
We can extend the definition of $`\mathrm{\Gamma }`$-convergence to families depending on a parameter $`\epsilon >0`$.
###### Definition 5.2.
For every $`\epsilon >0`$, let $`F_\epsilon `$ be a functional from $`U`$ to $`[\mathrm{},+\mathrm{}]`$. We say that $`F`$ is the $`\mathrm{\Gamma }(\tau )`$-limit of the family $`F_\epsilon `$, and write
$$F=\mathrm{\Gamma }(\tau )\text{-}\underset{\epsilon 0^+}{lim}F_\epsilon ,$$
if we have for every sequence $`\epsilon _h0^+`$
$$F=\mathrm{\Gamma }(\tau )\text{-}\underset{h+\mathrm{}}{lim}F_{\epsilon _h}.$$
Throughout this section, we work in the space $`L^p(\mathrm{\Omega },^m)`$ endowed with the strong topology. As pointed out in the introduction, we consider a non-negative function $`f=f(x,y^1,\mathrm{},y^n,z)`$ on $`\mathrm{\Omega }\times \mathrm{}^n\times 𝕄^{m\times d}`$ satisfying Assumptions 1, 2 and 3.
We fully characterize the $`\mathrm{\Gamma }(L^p)`$-limit of the family $`F_\epsilon :L^p(\mathrm{\Omega },^m)[0,+\mathrm{}]`$ where the functionals are defined by
$$F_\epsilon (u):=\{\begin{array}{cc}_\mathrm{\Omega }f(x,\frac{x}{\rho _1(\epsilon )},\mathrm{},\frac{x}{\rho _n(\epsilon )},u(x))𝑑x\hfill & \text{if}uW^{1,p}(\mathrm{\Omega },^m),\hfill \\ & \\ +\mathrm{}\hfill & \text{otherwise.}\hfill \end{array}$$
Precisely, this is our main result.
###### Theorem 5.3.
The family $`F_\epsilon `$$`\mathrm{\Gamma }(L^p)`$-converges and its $`\mathrm{\Gamma }(L^p)`$-limit $`F_{hom}:L^p(\mathrm{\Omega },^m)[0,+\mathrm{}]`$ is given by
$$F_{hom}(u)=\{\begin{array}{cc}_\mathrm{\Omega }f_{hom}(x,u(x))𝑑x\hfill & \text{if}uW^{1,p}(\mathrm{\Omega },^m),\hfill \\ & \\ +\mathrm{}\hfill & \text{otherwise},\hfill \end{array}$$
where $`f_{hom}`$ is obtained by the following cell problem
$$f_{hom}(x,z):=\underset{\varphi \mathrm{\Phi }}{inf}_\mathrm{}^nf(x,y,z+\underset{k=1}{\overset{n}{}}_{y^k}\varphi _k(y^1,\mathrm{},y^k))𝑑y$$
with the space $`\mathrm{\Phi }`$ defined by
$$\mathrm{\Phi }:=\underset{k=1}{\overset{n}{}}\mathrm{\Phi }_k\text{and}\mathrm{\Phi }_k:=L^p(\mathrm{}^{k1},W_{per}^{1,p}(\mathrm{},^m)).$$
###### Remark 5.4.
(i) Using the $`p`$-growth condition of $`f`$ and a density argument, it can be shown that
$$f_{hom}(x,z)=\underset{\varphi \mathrm{\Phi }_{reg}}{inf}_\mathrm{}^nf(x,y,z+\underset{k=1}{\overset{n}{}}_{y^k}\varphi _k(y^1,\mathrm{},y^k))𝑑y,$$
where
$$\mathrm{\Phi }_{reg}:=\underset{k=1}{\overset{n}{}}\mathrm{\Phi }_{k,reg}\text{and}\mathrm{\Phi }_{k,reg}:=C^1(\overline{\mathrm{}}^{k1},C_{per}^1(\mathrm{},^m)).$$
(ii) For every $`\delta >0`$ there exists a compact set $`X\mathrm{\Omega }`$ with $`\left|\mathrm{\Omega }\backslash X\right|\delta `$ such that the restriction of $`f`$ to $`X\times \mathrm{}^n\times 𝕄^{m\times d}`$ is continuous in $`(x,z)`$ for a.e. $`(y^1,\mathrm{},y^n)\mathrm{}^n`$ and so $`f_{hom}`$ is lower semicontinuous on $`X\times 𝕄^{m\times d}`$. In particular $`f_{hom}`$ is $`(\mathrm{\Omega })(𝕄^{m\times d})`$-measurable.
(iii) The convexity, the $`p`$-coerciveness and the $`p`$-growth condition on $`f`$ give the corresponding properties for the function $`f_{hom}`$. In particular $`F_{hom}`$ is continuous on $`W^{1,p}(\mathrm{\Omega },^m)`$, endowed with the strong topology.
Before proving the theorem, we state a series of lemmas. Only for simplicity of notations, we restrict ourselves to the case $`m=1`$. Fixed a sequence $`\epsilon _h0^+`$, we use for $`v_h`$ the same definition given in (3.1).
###### Lemma 5.5.
Let $`u_h`$ be a sequence converging weakly in $`W^{1,p}(\mathrm{\Omega })`$ to a function $`u`$. Then
$$\underset{h+\mathrm{}}{lim\; inf}_\mathrm{\Omega }f(x,v_h(x),u_h(x))𝑑x_\mathrm{\Omega }f_{hom}(x,u(x))𝑑x.$$
###### Proof.
For a suitable subsequence $`h_i`$ ,
$$\underset{i+\mathrm{}}{lim}_\mathrm{\Omega }f(x,v_{h_i}(x),u_{h_i}(x))𝑑x=\underset{h+\mathrm{}}{lim\; inf}_\mathrm{\Omega }f(x,v_h(x),u_h(x))𝑑x.$$
Refining the subsequence if necessary, we can suppose that $`u_{h_j}`$ generates a multiscale Young measure $`\nu `$. By Theorem 4.14 and Jensen’s inequality
$`\underset{i+\mathrm{}}{lim}{\displaystyle _\mathrm{\Omega }}f(x,v_{h_i}(x),u_{h_i}(x))𝑑x`$ $`{\displaystyle _{\mathrm{\Omega }\times \mathrm{}^n}}{\displaystyle _^d}f(x,y,z)𝑑\nu _{(x,y)}𝑑x𝑑y`$
$``$ $`{\displaystyle _{\mathrm{\Omega }\times \mathrm{}^n}}f(x,y,{\displaystyle _^d}z𝑑\nu _{(x,y)})𝑑x𝑑y`$
and by Theorems 3.6(ii) and 3.8
$``$ $`{\displaystyle _{\mathrm{\Omega }\times \mathrm{}^n}}f(x,y,u(x)+{\displaystyle \underset{k=1}{\overset{n}{}}}_{y^k}u_k(x,y^1,\mathrm{},y^k))𝑑x𝑑y`$
$``$ $`{\displaystyle _\mathrm{\Omega }}f_{hom}(x,u(x))𝑑x.`$
###### Lemma 5.6.
Let $`f:^d`$ be a convex function, such that for every $`z^d`$
$$\left|f(z)\right|c\left(b+\left|z\right|\right)^p,$$
(5.3)
where b and c are positive constants. Then, for all $`z_1,z_2^d`$
$$\left|f(z_1)f(z_2)\right|cd(1+2^p)\left(b+\left|z_1\right|+\left|z_2\right|\right)^{p1}\left|z_1z_2\right|.$$
(5.4)
The proof can be derived from \[16, Lemma 5.2\]. We observe that in (5.4) the estimate depends only by the costants b, c of growth condition (5.3) and not by the particular function $`f`$.
###### Lemma 5.7.
Let $`uW^{1,p}(\mathrm{\Omega })C^1(\mathrm{\Omega })`$. Then
$$\underset{u_hu}{inf}\left\{\underset{h+\mathrm{}}{lim\; sup}F_{\epsilon _h}(u_h)\right\}\underset{\psi \mathrm{\Psi }}{inf}_{\mathrm{\Omega }\times \mathrm{}^n}f(x,y,u(x)+\underset{k=1}{\overset{n}{}}_{y^k}\psi _k(x,y^1,\mathrm{},y^k))𝑑x𝑑y,$$
(5.5)
where the $`inf`$’s are made respectively on the sequences $`u_h`$ that converge strongly in $`L^p(\mathrm{\Omega })`$ to $`u`$ and on the space $`\mathrm{\Psi }`$ defined by
$$\mathrm{\Psi }:=\underset{k=1}{\overset{n}{}}\mathrm{\Psi }_k\text{and}\mathrm{\Psi }_k:=C^1(\overline{\mathrm{\Omega }}\times \overline{\mathrm{}}^{k1},C_{per}^1(\mathrm{})).$$
###### Proof.
Given an arbitrary function $`\psi =(\psi _1,\mathrm{},\psi _n)\mathrm{\Psi }`$, consider the sequence
$$u_h(x):=u(x)+\underset{k=1}{\overset{n}{}}\rho _k(\epsilon _h)\psi _k(x,v_h^k(x)),$$
where we used the short notation $`v_h^k(x):=(\frac{x}{\rho _1(\epsilon _h)},\mathrm{},\frac{x}{\rho _k(\epsilon _h)}).`$ We have $`u_hu`$ strongly in $`L^p(\mathrm{\Omega })`$ and $`u_h=u+_{k=1}^n_{y^k}\psi _k+r_h`$, with $`r_h0`$ strongly in $`L^p(\mathrm{\Omega },^d)`$.
The function $`g:\mathrm{\Omega }\times \mathrm{}^n`$ defined by
$$g(x,y):=f(x,y,u(x)+\underset{k=1}{\overset{n}{}}_{y^k}\psi _k(x,y^1,\mathrm{},y^k))$$
is admissible. Actually $`g𝒜dm^1`$, as evident by the estimate obtained through the $`p`$-growth condition:
$$\left|g(x,y)\right|c_2\left[1+(n+1)^{p1}\left(\left|u(x)\right|^p+\underset{k}{}M_k^p\right)\right],$$
where $`M_k:=sup_{\mathrm{\Omega }\times \mathrm{}^k}|_{y^k}\psi _k|`$.
By Lemma 5.6, the following inequality holds for some positive constants $`b`$, $`c`$ :
$$\left|g(x,v_h(x))f(x,v_h(x),u_h(x))\right|c\left|r_h(x)\right|\left(b+\left|u(x)\right|^{p1}+\left|r_h(x)\right|^{p1}\right).$$
By integrating over $`\mathrm{\Omega }`$, from Hölder’s inequality we obtain, for another positive constant $`c^{}`$,
$$_\mathrm{\Omega }|g(x,v_h(x))f(x,v_h(x),u_h(x))|dxc^{}_\mathrm{\Omega }|r_h(x)|^pdx$$
and thus Theorem 4.6 gives
$$\begin{array}{cc}\hfill \underset{h+\mathrm{}}{lim}& _\mathrm{\Omega }f(x,v_h(x),u_h(x))𝑑x=\underset{h+\mathrm{}}{lim}_\mathrm{\Omega }g(x,v_h(x))𝑑x\hfill \\ \hfill =& _{\mathrm{\Omega }\times \mathrm{}^n}g(x,y)𝑑x𝑑y=_{\mathrm{\Omega }\times \mathrm{}^n}f(x,y,u(x)+\underset{k=1}{\overset{n}{}}_{y^k}\psi _k(x,y^1,\mathrm{},y^k))𝑑x𝑑y.\hfill \end{array}$$
###### Definition 5.8.
We say that $`\mathrm{\Lambda }L^1(\mathrm{\Omega })`$ is an *inf-stable family* if, given $`\{\lambda _1,\mathrm{},\lambda _N\}\mathrm{\Lambda }`$ and $`\{\phi _1,\mathrm{},\phi _N\}C^1(\overline{\mathrm{\Omega }},[0,1])`$, with $`_{j=1}^N\phi _j=1`$ and $`N^+`$, there exists a $`\lambda \mathrm{\Lambda }`$ such that
$$\lambda \underset{j=1}{\overset{N}{}}\phi _j\lambda _j.$$
###### Lemma 5.9.
Let $`\mathrm{\Lambda }`$ be an inf-stable family of non-negative integrable functions on $`\mathrm{\Omega }`$. If for every $`\delta >0`$ there exists a compact set $`X_\delta \mathrm{\Omega }`$ such that $`\left|\mathrm{\Omega }\backslash X_\delta \right|\delta `$ and $`\lambda |_{X_\delta }`$ is continuous for each $`\lambda \mathrm{\Lambda }`$, then the function $`inf_{\lambda \mathrm{\Lambda }}\lambda `$ is measurable and the following commutation property holds:
$$\underset{\lambda \mathrm{\Lambda }}{inf}_\mathrm{\Omega }\lambda (x)𝑑x=_\mathrm{\Omega }\underset{\lambda \mathrm{\Lambda }}{inf}\lambda (x)dx.$$
(5.6)
This lemma can be derived by \[6, Lemma 4.3\] (see also ), by noting that for every $`\delta >0`$ $`inf_{\lambda \mathrm{\Lambda }}\lambda =\mathrm{ess}\mathrm{inf}_{\lambda \mathrm{\Lambda }}\lambda `$ on $`X_\delta `$. Anyway, we prefer to give a simple direct proof.
###### Proof.
Firstly we observe that for every $`\delta >0`$ the function $`inf_{\lambda \mathrm{\Lambda }}\lambda `$ is lower semicontinuous on $`X_\delta `$. In particular $`inf_{\lambda \mathrm{\Lambda }}\lambda `$ is measurable. By applying the Lindelöf theorem to each family $`\{E_\delta ^\lambda \}_{\lambda \mathrm{\Lambda }}`$, where
$$E_\delta ^\lambda :=\{(x,t)X_\delta \times :\lambda (x)<t\},$$
we can find a sequences $`\lambda _i`$ in $`\mathrm{\Lambda }`$ such that
$$\underset{\lambda \mathrm{\Lambda }}{inf}\lambda (x)=\underset{i}{inf}\lambda _i(x)\text{for }\text{a.e.}x\mathrm{\Omega }.$$
Fixed $`N^+`$ and $`\zeta >0`$, we choose a $`\delta >0`$ such that $`_{j=1}^N_{\mathrm{\Omega }\backslash X_\delta }\lambda _j\zeta `$. By the continuity property of the elements $`\lambda \mathrm{\Lambda }`$, the sets
$$A_i:=\{xX_\delta :\lambda _i(x)<\underset{1jN}{inf}\lambda _j(x)+\zeta \}$$
are open in $`X_\delta `$. Notice that $`X_\delta =_{i=1}^{\mathrm{}}A_i`$. For every $`i^+`$, let $`B_i`$ be a open subset of $`^d`$ for which $`B_iX_\delta =A_i`$ and let $`\{\phi _i\}_iC^1(\overline{\mathrm{\Omega }},[0,1])`$ be a partition of unity subordinate to $`\{B_i\}_i`$. By the inf-stability property, there exists a $`\lambda \mathrm{\Lambda }`$ such that $`\lambda _{j=1}^N\phi _j\lambda _j`$. We have
$$\begin{array}{cc}\hfill _\mathrm{\Omega }\lambda (x)𝑑x=& _{\mathrm{\Omega }\backslash X_\delta }\lambda (x)𝑑x+_{X_\delta }\lambda (x)𝑑x\hfill \\ \hfill & \underset{j=1}{\overset{N}{}}_{\mathrm{\Omega }\backslash X_\delta }\phi _j(x)\lambda _j(x)𝑑x+\underset{i=1}{\overset{\mathrm{}}{}}_{X_\delta }\phi _i(x)\lambda _i(x)𝑑x\hfill \\ \hfill & \zeta +_\mathrm{\Omega }\underset{1jN}{inf}\lambda _j(x)dx+\zeta \left|\mathrm{\Omega }\right|.\hfill \end{array}$$
Being $`N`$ and $`\zeta `$ arbitrary, the claim follows. ∎
We are now ready to assemble a proof of Theorem 5.3.
Let $`u_hu`$ in $`L^p(\mathrm{\Omega })`$. We want to show that $`lim\; infF_{\epsilon _h}(u_h)F_{hom}(u)`$. In this way inequality (5.1) will be proved. If $`lim\; infF_{\epsilon _h}(u_h)=+\mathrm{}`$, there is nothing to prove, so we can assume $`lim\; infF_{\epsilon _h}(u_h)<+\mathrm{}`$. For a suitable subsequence $`h_i`$ ,
$$\underset{i+\mathrm{}}{lim}F_{\epsilon _{h_i}}(u_{h_i})=\underset{h+\mathrm{}}{lim\; inf}F_{\epsilon _h}(u_h).$$
For $`i`$ large enough, $`F_{\epsilon _{h_i}}(u_{h_i})`$ is finite and therefore, by the definition of $`F_\epsilon `$, $`u_{h_i}W^{1,p}(\mathrm{\Omega })`$. Due to the $`p`$-coerciveness hypothesis on $`f`$, we can infer that $`u_{h_i}`$ is bounded in $`W^{1,p}(\mathrm{\Omega })`$. Refining the subsequence if necessary, we can suppose that $`u_{h_j}`$ converges weakly in $`W^{1,p}(\mathrm{\Omega })`$ to $`u`$ and thus we can apply Lemma 5.5.
It remains to check inequality (5.2). If $`uL^p(\mathrm{\Omega })W^{1,p}(\mathrm{\Omega })`$, then $`F_{hom}(u)=+\mathrm{}`$ and the inequality is obvious, while if $`uW^{1,p}(\mathrm{\Omega })`$, then we can apply Lemma 5.7 and, as in \[7, Theorem 3.3\], Lemma 5.9. In view of the density of $`W^{1,p}(\mathrm{\Omega })C^1(\mathrm{\Omega })`$ in $`W^{1,p}(\mathrm{\Omega })`$ and of the continuity of $`F_{hom}`$, by a standard diagonalization argument, it is not restrictive to assume that $`uW^{1,p}(\mathrm{\Omega })C^1(\mathrm{\Omega })`$.
For every $`\psi =(\psi _1,\mathrm{},\psi _n)\mathrm{\Psi }`$, define the function
$$\lambda _\psi (x):=_\mathrm{}^nf(x,y,u(x)+\underset{k=1}{\overset{n}{}}_{y^k}\psi _k(x,y^1,\mathrm{},y^k))𝑑y.$$
We claim that the family $`\mathrm{\Lambda }:=\{\lambda _\psi :\psi \mathrm{\Psi }\}`$ satisfies the hypotheses of Lemma 5.9. In fact, from the $`p`$-growth condition on $`f`$, it is easy to show that each function in $`\mathrm{\Lambda }`$ is integrable on $`\mathrm{\Omega }`$. Moreover, by Remark 5.4(ii), for every $`\delta >0`$ there exists a compact set $`X\mathrm{\Omega }`$ with $`\left|\mathrm{\Omega }\backslash X\right|\delta `$ such that $`\lambda _\psi `$ is continuous on $`X`$ for each $`\psi \mathrm{\Psi }`$. It remains to prove the inf-stability.
Given $`\{\psi ^{(1)},\mathrm{},\psi ^{(N)}\}\mathrm{\Psi }`$ and $`\{\phi _1,\mathrm{},\phi _N\}C^1(\overline{\mathrm{\Omega }},[0,1])`$, with $`_{j=1}^N\phi _j=1`$ and $`N^+`$, consider the function
$$\psi :=(\underset{j=1}{\overset{N}{}}\phi _j\psi _1^{(j)},\mathrm{},\underset{j=1}{\overset{N}{}}\phi _j\psi _n^{(j)})\mathrm{\Psi }.$$
Thanks to the convexity of $`f`$, we have $`\lambda _\psi _{j=1}^N\phi _j\lambda _{\psi ^{(j)}}`$:
$$\begin{array}{cc}\hfill \lambda _\psi (x)=& _\mathrm{}^nf(x,y,u(x)+\underset{j=1}{\overset{N}{}}\underset{k=1}{\overset{n}{}}_{y^k}\phi _j(x)\psi _k^{(j)}(x,y^1,\mathrm{},y^k))𝑑y\hfill \\ \hfill =& _\mathrm{}^nf(x,y,\underset{j=1}{\overset{N}{}}\phi _j(x)\left(u(x)+\underset{k=1}{\overset{n}{}}_{y^k}\psi _k^{(j)}(x,y^1,\mathrm{},y^k)\right))𝑑y\hfill \\ \hfill & \underset{j=1}{\overset{N}{}}\phi _j(x)_\mathrm{}^nf(x,y,u(x)+\underset{k=1}{\overset{n}{}}_{y^k}\psi _k^{(j)}(x,y^1,\mathrm{},y^k))𝑑y=\underset{j=1}{\overset{N}{}}\phi _j(x)\lambda _{\psi ^{(j)}}(x).\hfill \end{array}$$
Finally, by inequality (5.5), equality (5.6) and Remark 5.4(i),
$`\underset{u_hu}{inf}`$ $`\left\{\underset{h+\mathrm{}}{lim\; sup}F_{\epsilon _h}(u_h)\right\}{\displaystyle _\mathrm{\Omega }}\underset{\psi \mathrm{\Psi }}{inf}\left({\displaystyle _\mathrm{}^n}f(x,y,u(x)+{\displaystyle \underset{k=1}{\overset{n}{}}}_{y^k}\psi _k(x,y^1,\mathrm{},y^k))𝑑y\right)dx`$
$`{\displaystyle _\mathrm{\Omega }}\underset{\varphi \mathrm{\Phi }_{reg}}{inf}\left({\displaystyle _\mathrm{}^n}f(x,y,u(x)+{\displaystyle \underset{k=1}{\overset{n}{}}}_{y^k}\varphi _k(y^1,\mathrm{},y^k))𝑑y\right)dx={\displaystyle _\mathrm{\Omega }}f_{hom}(x,u(x))𝑑x.`$
The proof is complete.
$`\mathrm{}`$
Assumpion 3 cannot be weakened too much: even if $`f`$ is a Borel function, $`\mathrm{\Gamma }`$-convergence Theorem 5.3 may be not applicable, as shown in the following examples (see also \[12, Example 3.1\]).
###### Example 5.10.
Let $`A_i`$ be the Borel sets defined as in Example 4.8 and let $`B:=_{i=1}^{\mathrm{}}A_{2i}`$. Notice that $`_{ij}(A_iA_j)`$ is countable. In the case $`d=n=m=1`$, $`\rho _1(\epsilon _h)=h^1`$ and $`\mathrm{\Omega }=(0,1)`$, consider the Borel function $`f(x,y,z):=[2\chi _B(x,y)]\left|z\right|^p`$. We remark that $`f`$ satisfies only Assumptions 1 and 2. We have for every $`uW^{1,p}((0,1))`$
$$F_h(u)=_0^1f(x,hx,u(x))𝑑x=\{\begin{array}{cc}_0^1\left|u(x)\right|^p𝑑x\hfill & \text{if}h0mod2,\hfill \\ & \\ 2_0^1\left|u(x)\right|^p𝑑x\hfill & \text{if}h1mod2.\hfill \end{array}$$
Clearly the sequence $`F_h`$ is not $`\mathrm{\Gamma }`$-convergent in $`L^p((0,1))`$ with respect to the strong topology.
###### Example 5.11.
Let $`d=m=1`$, $`n=2`$, $`\rho _1(\epsilon _h)=h^1`$, $`\rho _2(\epsilon _h)=h^2`$ and $`\mathrm{\Omega }=(0,1)`$. Consider the Borel function $`f(x,y^1,y^2,z):=[2\chi _B(y^1,y^2)]\left|z\right|^p`$, where $`B`$ is defined as in the former example. Even if $`f`$ does not depend by $`x`$ and satisfies Assumptions 1 and 2, the sequence $`F_h`$ is not $`\mathrm{\Gamma }(L^p)`$-convergent.
## 6. Iterated homogenization
The homogenized function $`f_{hom}`$ can be obtained also by the following iteration:
$$\begin{array}{cc}& f_{hom}^{[n]}(x,y^1,\mathrm{},y^{n1},z):=\underset{\varphi W_{per}^{1,p}(\mathrm{},^m)}{inf}_{\mathrm{}}f(x,y^1,\mathrm{},y^n,z+\varphi (y^n))𝑑y^n,\hfill \\ & f_{hom}^{[n1]}(x,y^1,\mathrm{},y^{n2},z):=\underset{\varphi W_{per}^{1,p}(\mathrm{},^m)}{inf}_{\mathrm{}}f_{hom}^{[n]}(x,y^1,\mathrm{},y^{n1},z+\varphi (y^{n1}))𝑑y^{n1},\hfill \\ & \mathrm{}\hfill \\ & f_{hom}(x,z)=f_{hom}^{[1]}(x,z):=\underset{\varphi W_{per}^{1,p}(\mathrm{},^m)}{inf}_{\mathrm{}}f_{hom}^{[2]}(x,y^1,z+\varphi (y^1))𝑑y^1.\hfill \end{array}$$
###### Remark 6.1.
(i)The convexity, the $`p`$-coerciveness and the $`p`$-growth condition on $`f`$ give the corresponding properties for the function $`f_{hom}^{[n]}`$. Moreover, $`f_{hom}^{[n]}`$ is still an admissible integrand. In fact, for every $`\delta >0`$ there exist a compact set $`X\mathrm{\Omega }`$ with $`\left|\mathrm{\Omega }\backslash X\right|\delta `$ and a compact set $`Y\mathrm{}`$ with $`\left|\mathrm{}\backslash Y\right|\delta `$, such that the restriction of $`f`$ to $`X\times Y^{n1}\times \mathrm{}\times 𝕄^{m\times d}`$ is continuous in $`(x,y^1,\mathrm{},y^{n1},z)`$ for a.e. $`y^n\mathrm{}`$. Consequently, following closely \[15, Lemma 4.1\], it can be proved that $`f_{hom}^{[n]}`$ is continuous on $`X\times Y^{n1}\times 𝕄^{m\times d}`$.
(ii)Clearly, the properties of $`f_{hom}^{[n]}`$ give the corresponding ones for $`f_{hom}^{[n1]}`$ and so on.
We prove only the inequality $`f_{hom}f_{hom}^{[1]}`$, since the opposite inequality comes directly. Fixed $`(x,z)\mathrm{\Omega }\times 𝕄^{m\times d}`$ and $`\varphi _k\mathrm{\Phi }_{k,reg}`$ for $`k=1,\mathrm{},n1`$, by using a commutation argument as Lemma 5.9, we get
$`\underset{\varphi _n\mathrm{\Phi }_{n,reg}}{inf}`$ $`{\displaystyle _\mathrm{}^n}f(x,y,z+{\displaystyle \underset{k=1}{\overset{n}{}}}_{y^k}\varphi _k(y^1,\mathrm{},y^k))𝑑y`$
$`=`$ $`{\displaystyle _{\mathrm{}^{n1}}}\underset{\varphi _n\mathrm{\Phi }_{n,reg}}{inf}\left({\displaystyle _{\mathrm{}}}f(x,y,z+{\displaystyle \underset{k=1}{\overset{n}{}}}_{y^k}\varphi _k(y^1,\mathrm{},y^k))𝑑y^n\right)dy^1\mathrm{}dy^{n1}`$
$``$ $`{\displaystyle _{\mathrm{}^{n1}}}f_{hom}^{[n]}(x,y^1,\mathrm{},y^{n1},z+{\displaystyle \underset{k=1}{\overset{n1}{}}}_{y^k}\varphi _k(y^1,\mathrm{},y^k))𝑑y^1\mathrm{}𝑑y^{n1}.`$
By repeating the commutation procedure, we obtain
$`\underset{\varphi _{n1}\mathrm{\Phi }_{n1,reg}}{inf}`$ $`{\displaystyle _{\mathrm{}^{n1}}}f_{hom}^{[n]}(x,y^1,\mathrm{},y^{n1},z+{\displaystyle \underset{k=1}{\overset{n1}{}}}_{y^k}\varphi _k(y^1,\mathrm{},y^k))𝑑y^1\mathrm{}𝑑y^{n1}`$
$``$ $`{\displaystyle _{\mathrm{}^{n2}}}f_{hom}^{[n1]}(x,y^1,\mathrm{},y^{n2},z+{\displaystyle \underset{k=1}{\overset{n2}{}}}_{y^k}\varphi _k(y^1,\mathrm{},y^k))𝑑y^1\mathrm{}𝑑y^{n2}`$
and so on. Then
$$f_{hom}(x,z)\underset{\varphi _1\mathrm{\Phi }_{1,reg}}{inf}\mathrm{}\underset{\varphi _n\mathrm{\Phi }_{n,reg}}{inf}_\mathrm{}^nf(x,y,z+\underset{k=1}{\overset{n}{}}_{y^k}\varphi _k(y^1,\mathrm{},y^k))𝑑yf_{hom}^{[1]}(x,z).$$
Acknowledgments
I wish to thank Gianni Dal Maso for many helpful and interesting discussions. |
warning/0506/astro-ph0506760.html | ar5iv | text | # Near Infrared Survey of Populous Clusters in the LMC: Preliminary Results
## 1. Introduction
The Large Magellanic Cloud (LMC) is an attractive astronomical target for a variety of reasons. Owing to its relative proximity, stellar populations in the LMC can be easily resolved. These populations exhibit an array of star formation processes and episodes in a dynamic environment making the LMC well suited for studying the formation and evolution of a satellite galaxy. Traditionally, the LMC has been thought of as an approximately planar galaxy that, in spite of its proximity, can be assumed to lie at a single distance from us. In contrast, Caldwell & Coulson (1986) have shown that the LMC disk is tilted with respect to the plane of the sky. More recent work has not only confirmed that the LMC is tilted, but it also indicates that the LMC disk is considerably thicker than previously assumed.
van der Marel & Cioni (2001) determined, through the use of red giant branch and asymptotic giant branch stars as relative distance indicators, that the LMC is tilted $`34.^{}7\pm 6.^{}2`$ with respect to the line of sight ($`0^{}`$ is face-on) such that the Northeast portion of the LMC is closer to us than the Southwest. Olsen & Salyk (2002) confirmed this result, finding $`i=35.^{}8\pm 2.^{}4`$ by utilizing core He burning red clump (RC) stars as their relative distance indicator. Additionally, by studying carbon star kinematics in the LMC disk, van der Marel et al. (2002) have determined that $`v/\sigma =2.9\pm 0.9`$, implying that the LMC disk is thicker than the Milky Way thick disk ($`v/\sigma 3.9`$). Lastly, the Magellanic Stream (e.g. Putman et al. 2003), flaring (Alves & Nelson, 2000) and elongation of the LMC disk (van der Marel & Cioni 2001) and the possibility that the LMC disk is warped (Olsen & Salyk 2002, Nikolaev et al. 2004) all indicate that the LMC has not escaped unharmed from its tidal interactions with the Milky Way and Small Magellanic Cloud.
The distance to the LMC has been a topic of considerable discussion in recent years and a variety of methods have been employed to calculate this distance; e.g. variable stars (Cepheids, RR Lyraes, Miras), color-magnitude diagram (CMD) features (main sequence turn off, tip of the red giant branch, RC stars), and SN 1987a. There has been, until recently, little agreement between the different methods and sometimes even amongst distances calculated using a single method. This lead to a “long” and “short” distance scale for the LMC with a “short” distance modulus, $`(mM)_0`$, of $``$18.2-18.3 mag and a “long” distance of $``$18.5-18.7 mag. Clementini et al. (2003) demonstrate this distance problem (top panel, their Fig. 8), and find that the long and short distance scale can be reconciled, at least to within the errors, with improved photometry and/or reddening estimates (bottom panel, their Fig. 8) for some of the previous works.
A primary reason for interest in the LMC distance is its use as the extragalactic distance scale zeropoint. The Hubble Space Telescope Key Project to determine $`H_0`$ (Freedman et al. 2001) utilized a sample of LMC Cepheid variables to define the fiducial period-luminosity relation. Cepheid distances were then used to calibrate secondary standard candles which lie further along the extragalactic distance ladder. Thus, the accuracy of their determination of $`H_0`$ ($`72\pm 8`$ km s<sup>-1</sup> Mpc<sup>-1</sup>) hinges on the accuracy of their zeropoint, $`(mM)_{0,LMC}=18.5\pm 0.10`$. It turns out the error in their calculation is dominated by the uncertainty in $`(mM)_{0,LMC}`$; it takes up $`6.5\%`$ of their $`9\%`$ error budget (Mould et al. 2000).
In this paper we will present preliminary results from our near-infrared survey of populous clusters in the LMC. Section 2 presents our observations of LMC cluster and field stars. In the next two sections we discuss our application of the $`K`$-band luminosity of the RC as a standard candle for calculating absolute cluster distances (§3) and for determining relative distances to the LMC fields (§4). Finally, in §5 we talk about our future work on this project.
## 2. Data
### 2.1. Observations
We have obtained near infrared images for a sample of intermediate age LMC clusters over the course of two, three night observing runs (20-22 January 2003 and 06-08 February 2004) at the CTIO 4m. Our observations were made with the Infrared Side Port Imager (ISPI) which utilizes a 2048 $`\times `$ 2048 pixel HAWAII 2 HgCdTe array. In the f/8 configuration, ISPI yields an 11 $`\times `$ 11 field of view and a plate scale of $``$ 0$`.^{\prime \prime }`$33 pixel<sup>-1</sup>. For our observations we used a nine-point dither pattern, centered on each cluster, and total integration times as follows: J = 540s, H = 846s, K = 846s. Average seeing for all six nights was $``$ 1.2<sup>′′</sup>. Table 1 lists two of our 18 program clusters along with right ascention and declination (J2000) and passbands in which the clusters were observed.
### 2.2. Reduction
All data was processed using standard data reduction steps, which we will now summarize. Images were dark subtracted, sky subtracted and then flat fielded using on-off dome flats. Due to the combination of ISPI’s wide field of view and the relatively large steps in our dither pattern ($``$ 30), these images suffer from geometric distortions, caused mostly by the curvature of the focal plane. To correct for this, we apply a high order distortion correction to each image using the IRAF task GEOTRAN. The corrected images are then aligned, shifted and averaged to create a final science frame for each cluster and filter.
Science frames were photometered using a combination of DAOPHOT and ALLSTAR (Stetson 1987) as follows. A rough PSF was constructed using the brightest $``$ 200 stars in each image. The rough PSF was then used to remove neighbors from around the PSF stars, allowing the creation of a more robust PSF from the cleaned image. ALLSTAR was utilized to fit this improved PSF to all stars in the science frame. In an effort to detect and photometer faint stars and/or companions, we used a single iteration of subtracting all stars detected and fit in the first ALLSTAR pass, then searching for previously undetected stars in our fields. All new detections were run through ALLSTAR using the same PSF as in the first ALLSTAR pass and these stars were added to the photometry list. At this point, aperture corrections, calculated for each frame, were applied to the PSF photometry. Finally, aperture corrected photometry lists from each filter were combined with the requirement that a star be detected in all filters for it to be kept in the final combined list of instrumental magnitudes. Zero points and color tranformations appropriate for our data were calculated by comparing our instrumental magnitudes with photometry from the 2MASS All-Sky Data Release<sup>1</sup><sup>1</sup>1http://www.ipac.caltech.edu/2mass/releases/allsky for each field in our program.
## 3. Preliminary Distances
Figure 1 presents ($`K`$, $`JK`$) CMDs for NGC 1651 and Hodge 4, where all stars within $`1^{}`$ of the cluster centers are shown. Both CMDs show a prominent RC at $`K16.9`$ and a well populated RGB extending up to $`K12.5`$.
We follow the method of Grocholski & Sarajedini (2002) in using a box that extends 0.8 mag in $`K`$ and 0.2 mag in $`JK`$ (shown in Fig. 1), centered by eye, to select the RC stars. $`K_{RC}`$ is calculated by taking the median value of all stars within this box. For NGC 1651, $`K_{RC}=16.93\pm 0.02`$ and $`K_{RC}=16.81\pm 0.02`$ for Hodge 4.
With regards to the reddening of each cluster, we utilize the dust maps of Burstein & Heiles (1982) and Schlegel, Finkbeiner, & Davis (1998). Since the values determined from both dust maps, for each cluster, are in good agreement, we adopt the average value from the two maps as our cluster reddenings. For NGC 1651 we find $`E(BV)=0.12\pm 0.02`$ and for Hodge 4, $`E(BV)=0.05\pm 0.01`$. Using the relations from Cardelli, Clayton, & Mathis (1989), $`A_V=3.1E(BV)`$ and $`A_K=0.11A_V`$, these reddenings translate to $`A_K=0.041\pm 0.003`$ and $`A_K=0.017\pm 0.007`$.
Previous authors have shown that the absolute RC magnitude varies as a function of age and metallicity for visible and near-infrared bands, with this variation seen in both theoretical (e.g. Salaris & Girardi 2002) and observational data (e.g. Sarajedini 1999, Cole 1998). An in-depth comparison of the absolute $`K`$-band RC magnitude with age and metallicity for a sample of simple stellar populations was performed by Grocholski & Sarajedini (2002). These authors advocate using an interpolation over either their observational data or the theoretical models of Girardi & Salaris (2001) to create an $`M_K(RC)`$ “plane” which, given a cluster’s age and metallicity, can be used to predict $`M_K(RC)`$ for that cluster.
In many cases, however, age and metallicity values for LMC clusters are not readily available or not reliable. For example, Olszewski et al. (1991) have presented the only large scale determination of metallicities for LMC clusters, based on the Ca II triplet. However, many of their cluster \[Fe/H\] values, including NGC 1651 and Hodge 4, are based on observations of a single star. Sarajedini et al. (2002) calculated \[Fe/H\] values for both of these clusters using the slope of the RGB. While their value for Hodge 4 is consistent with that of Olszewski et al. (1991), their value for NGC 1651 is 0.3 dex more metal rich. As such, in the current work we choose not to apply the full calibration of $`M_K(RC)`$ as discussed above, but rather we adopt the value of $`M_K(RC)=1.61\pm 0.04`$ given in Grocholski & Sarajedini (2002). We note that Sarajedini et al. (2002), utilizing the full RC treatment from Grocholski & Sarajedini (2002), find $`M_K(RC)=1.56\pm 0.12`$ for NGC 1651 and $`M_K(RC)=1.64\pm 0.17`$. This implies that error in the value we have chosen to use for $`M_K(RC)`$ is likely larger than that quoted.
Using the values listed above for $`K(RC)`$, $`M_K(RC)`$, and $`A_K`$ for each custer, we find for NGC 1651, $`(mM)_0=18.50\pm 0.06`$ and $`(mM)_0=18.40\pm 0.05`$ for Hodge 4. The errors quoted are the random errors added in quadrature. These numbers are consistent with the LMC distance, $`(mM)_0=18.50\pm 0.10`$ used in the HST Key Project to determine an accurate value of $`H_0`$ (see Freedman et al. 2001 for more information). Additionally, these distances agree with the LMC geometry determined by van der Marel & Cioni (2001) and Olsen & Salyk (2002) in that Hodge 4 should be closer to us than NGC 1651, based on the tilt of the LMC’s disk and the location of the clusters in the LMC. Lastly, Sarajedini et al. (2002) find, for NGC 1651 and Hodge 4, $`(mM)_0=18.55\pm 0.12`$ and $`18.52\pm 0.17`$. These distances are in agreement, within the errors, with our results.
## 4. Field Stars
In Figure 2 we show ($`K`$, $`JK`$) CMDs for the field stars surrounding our clusters and, as with the clusters, the RC and RGB for the two fields are easily visible. The major difference between the cluster and field CMDs is the wider field RGB (spread in color) and larger field RC (spread in both color and luminosity) caused by the intrinsic distribution in age and metallicity of the field population. Although this situation is a bit more complicated for the RC (see Salaris & Girardi 2002; Grocholski & Sarajedini 2002), in general, older and/or more metal rich populations have redder RGBs and redder and fainter RCs than young and/or metal poor populations. Due to the mixed population in the LMC field, dealing with the field RC luminosity as a standard candle becomes a much more formidable task than dealing with a simple stellar population RC.
However, if we can make the assumption that each observed field within the LMC has a similar mix of stars in terms of age and metallicity, then the field RC can be easily applied as a relative distance indicator since $`M_K(RC)`$ should be the same for all fields. This assumption is reasonable given that the bulk of the LMC RC stars are $`4`$ Gyr old (Girardi & Salaris 2001) and differential rotation should destroy any record of the initial local age and metallicity distribution on timescales much shorter than this (Olsen and Salyk 2002). Additionally, if this assumption holds true, the variation in RC colors amongst our fields is indicative of the relative reddenings of these fields.
Similar to §3, we determine the magnitude and color values of the field RC by taking the median value of all stars within a predefined box, centered on the RC. Since the field RC is more extended than the cluster RCs, we have doubled the size of our box to 1.6 mag in $`K`$ and 0.4 mag in $`(JK)`$, as shown in Fig. 2. For our NGC 1651 and Hodge 4 fields, respectively, we find $`K_{RC}=17.01\pm 0.01`$, $`(JK)_{RC}=0.630\pm 0.002`$ and $`K_{RC}=16.85\pm 0.01`$, $`(JK)_{RC}=0.662\pm 0.002`$. Applying the reddening corrections we derived from the dust maps (section 3), we find $`K_0(RC)=16.97\pm 0.01`$ and $`K_0(RC)=16.83\pm 0.01`$ for the NGC 1651 and Hodge 4 fields. This result implies that the field around NGC 1651 is $`0.14`$ mag farther from us than the Hodge 4 field, consistent with the results from §3. However, assuming the difference in color is solely due to a change in the reddening between these two fields, the Hodge 4 field suffers from $`E(BV)=0.061`$ more extinction than the NGC 1651 field, in the opposite sense of what is predicted by both dust maps discussed in §3. This could indicate that the field populations are in fact not composed of the same stellar mixture (the Hodge 4 field may be older and/or more metal rich), which would render this relative distance method invalid for our data. It is also possible that there exists a problem in our photometric calibration, leading to the discrepency between our results and the reddening maps of Schlegel et al. (1998) and Burstein & Heiles (1982). We are currently exploring the cause of this disagreement.
## 5. Future Work
As mentioned previously, the two clusters presented here are only a sample of our entire data set. We plan to use these data to address two problems. First, we will utilize $`K(RC)`$ values for each cluster, along with the full method of determining $`M_K(RC)`$, to calculate the individual cluster distances. (We note that we are in the process of measuring homogeneous ages and metallicities for each of our clusters, which will be used in properly calculating $`M_K(RC)`$) This will allow us to determine an accurate distance for the LMC. Additionally, with the areal coverage of our clusters, we will be in a position to compare the distribution of our program clusters with the LMC geometry derived by van der Marel and Cioni (2001). Second, the plethora of field stars surrounding each of our clusters provides an opportunity to study the field populations in the LMC. If the assumption that RC stars are evenly mixed throughout the LMC is confirmed, then we will be able to determine the relative distance to each of these fields, thereby studying the three dimensional distribution of the LMC disk. This information will provide insight into the reality of the “warp” found in the southwest portion of the LMC disk (Olsen and Salyk 2002). We will also be in a position to compare the relative distributions of our LMC clusters with their surrounding fields and explore whether or not the LMC clusters occupy the same plane as the LMC disk.
#### Acknowledgments.
This research is supported by NSF CAREER grant AST-0094048 to A. Sarajedini. |
warning/0506/math0506400.html | ar5iv | text | # An infinite genus mapping class group and stable cohomology11footnote 1 This version: January 30, 2008. L.F. was partially supported by the ANR Repsurf:ANR-06-BLAN-0311. Available electronically at http://www-fourier.ujf-grenoble.fr/~funar
## 1 Introduction
### 1.1 Statements of the main results
The tower of all extended mapping class groups was considered first by Moore and Seiberg () as part of the conformal field theory data. This object is actually a groupoid, which has been proved to be finitely presented (see ). When seeking for a group analog Penner () investigated a universal mapping class group which arises by means of a completion process and which is closely related to the group of homeomorphisms of the circle, but it seems to be infinitely generated.
In , we introduced the universal mapping class group in genus zero $``$. The latter is an extension of the Thompson’s group $`V`$ (see ) by the infinite spherical pure mapping class group. We proved in that the group $``$ is finitely presented and we exhibited an explicit presentation. Our main difference with the previous attempts is that we consider groups acting on infinite surfaces with a prescribed behaviour at infinity that comes from actions on trees.
Following the same kind of approach, we propose a treatment of the arbitrary genus case by introducing a mapping class group $``$, called the asymptotic infinite genus mapping class group, that contains a large part of the mapping class groups of compact surfaces with boundary. More precisely, the group $``$ contains all the pure mapping class groups $`P(\mathrm{\Sigma }_{g,n})`$ of compact surfaces $`\mathrm{\Sigma }_{g,n}`$ of genus $`g`$ with $`n`$ boundary components, for any $`g0`$ and $`n>0`$. Its construction is roughly as follows. Let $`𝒮`$ denote the surface obtained by taking the boundary of the 3-dimensional thickening of the complete trivalent tree, and further let $`𝒮_{\mathrm{}}`$ be the result of attaching a handle to each cylinder in $`𝒮`$ that corresponds to an edge of the tree (see figure 1). Then $``$ is the group of mapping classes of those homeomorphisms of $`𝒮_{\mathrm{}}`$ which preserve a certain rigid structure at infinity (see Definition 1.3 for the precise definition). This rigidity condition essentially implies that $``$ induces a group of transformations on the set of ends of the tree, which is isomorphic to Thompson’s group $`V`$. The relation between both groups is enlightened by a short exact sequence $`1PV1`$, where $`P`$ is the mapping class group of compactly supported homeomorphisms of $`𝒮_{\mathrm{}}`$. The latter is an infinitely generated group. Our first result is:
###### Theorem 1.1.
The group $``$ is finitely generated.
The interest in considering the group $``$, outside the framework of the topological quantum field theory where it can replace the duality groupoid, is the following homological property:
###### Theorem 1.2.
The rational homology of $``$ is isomorphic to the stable rational homology of the (pure) mapping class groups.
As a corollary of the argument of the proof (see Proposition 3.1), the group $``$ is perfect, and $`H_2(,)=`$. For a reason that will become clear in what follows, the generator of $`H^2(,)`$ is called the first universal Chern class of $``$, and is denoted $`c_1()`$.
Let $`_g`$ be the mapping class group of a closed surface $`\mathrm{\Sigma }_g`$ of genus $`g`$. We show that the standard representation $`\rho _g:_g\mathrm{Sp}(2g,)`$ in the symplectic group, deduced from the action of $`_g`$ on $`H_1(\mathrm{\Sigma }_g,)`$, extends to the infinite genus case, by replacing the finite dimensional setting by concepts of Hilbertian analysis. In particular, a key role is played by Shale’s restricted symplectic group $`\mathrm{Sp}_{\mathrm{res}}(_r)`$ on the real Hilbert space $`_r`$ generated by the homology classes of non-separating closed curves of $`𝒮_{\mathrm{}}`$. We have then:
###### Theorem 1.3.
The action of $``$ on $`H_1(𝒮_{\mathrm{}},)`$ induces a representation $`\rho :\mathrm{Sp}_{\mathrm{res}}(_r)`$.
The generator $`c_1`$ of $`H^2(_g,)`$ is called the first Chern class, since it may be obtained as follows (see, e.g., ). The group $`\mathrm{Sp}(2g,)`$ is contained in the symplectic group $`\mathrm{Sp}(2g,)`$, whose maximal compact subgroup is the unitary group $`U(g)`$. Thus, the first Chern class may be viewed in $`H^2(B\mathrm{Sp}(2g,),)`$. It can be first pulled-back on $`H^2(B\mathrm{Sp}(2g,)^\delta ,)=H^2(\mathrm{Sp}(2g,),)`$ and then on $`H^2(_g,)`$ via $`\rho _g`$. This is the generator of $`H^2(_g,)`$. Here $`B\mathrm{Sp}(2g,)^\delta `$ denotes the classifying space of the group $`\mathrm{Sp}(2g,)`$ endowed with the discrete topology.
The restricted symplectic group $`\mathrm{Sp}_{\mathrm{res}}(_r)`$ has a well-known 2-cocycle, which measures the projectivity of the Berezin-Segal-Shale-Weil metaplectic representation in the bosonic Fock space (see , Chapter 6 and Notes p. 171). Contrary to the finite dimension case, this cocycle is not directly related to the topology of $`\mathrm{Sp}_{\mathrm{res}}(_r)`$, since the latter is a contractible Banach-Lie group. However, $`\mathrm{Sp}_{\mathrm{res}}(_r)`$ embeds into the restricted linear group of Pressley-Segal $`\mathrm{GL}_{\mathrm{res}}^0()`$ (see ), where $``$ is the complexification of $`_r`$, which possesses a cohomology class of degree 2: the Pressley-Segal class $`PSH^2(\mathrm{GL}_{\mathrm{res}}^0(),^{})`$. The group $`\mathrm{GL}_{\mathrm{res}}^0()`$ is a homotopic model of the classifying space $`BU`$, where $`U=\underset{n\mathrm{}}{lim}U(n,)`$, and the class $`PS`$ does correspond to the universal first Chern class. Its restriction on $`\mathrm{Sp}_{\mathrm{res}}(_r)`$ is closely related to the Berezin-Segal-Shale-Weil cocycle, and reveals the topological origin of the latter. Via the composition of morphisms
$$\mathrm{Sp}_{\mathrm{res}}(_r)\mathrm{GL}_{\mathrm{res}}^0(),$$
we then derive from $`PS`$ an integral cohomology class on $``$ (see Theorem 5.1 for a more precise statement):
###### Theorem 1.4.
The Pressley-Segal class $`PSH^2(\mathrm{GL}_{\mathrm{res}}^0(),^{})`$ induces the first universal Chern class $`c_1()H^2(,)`$.
Acknowledgements. The authors are indebted to Vlad Sergiescu for enlighting discussions and particularly for suggesting the existence of a connection between the first universal Chern class of $``$ and the Pressley-Segal class. They are thankful to the referees for suggestions improving the exposition.
### 1.2 Definitions
#### 1.2.1 The infinite genus mapping class group $``$
Set $`(\mathrm{\Sigma }_{g,n})`$ for the extended mapping class group of the $`n`$-holed orientable surface $`\mathrm{\Sigma }_{g,n}`$ of genus $`g`$, consisting of the isotopy classes of orientation-preserving homeomorphisms of $`\mathrm{\Sigma }_{g,n}`$ which respect a fixed parametrization of the boundary circles, allowing them to be permuted among themselves.
We wish to construct a mapping class group, containing all mapping class groups $`(\mathrm{\Sigma }_{g,n})`$. It seems impossible to construct such a group, but if one relaxes slightly our requirements then we could follow our previous method used for the genus zero case in .
The choice of the extra structure involved in the definitions below is important because the final result might depend on it. For instance, using the same planar punctured surface but different decompositions one obtained in two non-isomorphic braided Ptolemy-Thompson groups.
###### Definition 1.1 (The infinite genus surface $`𝒮_{\mathrm{}}`$).
Let $`𝒯`$ be the complete trivalent planar tree and $`𝒮`$ be the surface obtained by taking the boundary of the 3-dimensional thickening of $`𝒯`$.
By grafting an edge-loop (i.e. the graph obtained by attaching a loop to a boundary vertex of an edge) at the midpoint of each edge of $`𝒯`$, one obtains the graph $`𝒯_{\mathrm{}}`$. The surface $`𝒮_{\mathrm{}}`$ is the boundary of the 3-dimensional thickening of $`𝒯_{\mathrm{}}`$.
The graph $`𝒯`$ (respectively $`𝒯_{\mathrm{}}`$) is embedded in $`𝒮`$ (respectively $`𝒮_{\mathrm{}}`$) as a cross-section of the fiber projection, as indicated on figure 1. Thus, $`𝒮_{\mathrm{}}`$ is obtained by removing small disks from $`𝒮`$ centered at midpoints of edges of $`𝒯`$ and gluing back one holed tori $`\mathrm{\Sigma }_{1,1}`$, called wrists which correspond to the thickening of edge-loops.
It is convenient to assume that $`𝒯`$ is embedded in a horizontal plane, while the edge-loops are in vertical planes (see figure 1).
###### Definition 1.2 (Pants decomposition of $`𝒮_{\mathrm{}}`$).
A pants decomposition of the surface $`𝒮_{\mathrm{}}`$ is a maximal collection of distinct nontrivial simple closed curves on $`𝒮_{\mathrm{}}`$ which are pairwise disjoint and non-isotopic. The complementary regions (which are 3-holed spheres) are called pairs of pants.
By construction, $`𝒮_{\mathrm{}}`$ is naturally equipped with a pants decomposition, which will be referred to below as the canonical (pants) decomposition, as shown in figure 1:
* the wrists are decomposed using a meridian circle and the boundary circle of $`\mathrm{\Sigma }_{1,1}`$.
* there is one pair of pants for each edge, which has one boundary circle for attaching the wrists, and two circles to grip to the other type of pants. We call them edge pants.
* there is one pair of pants for each vertex of the tree, called vertex pants.
A pants decomposition is asymptotically trivial if outside a compact subsurface of $`𝒮_{\mathrm{}}`$, it coincides with the canonical pants decomposition.
###### Definition 1.3.
1. A connected subsurface $`\mathrm{\Sigma }`$ of $`𝒮_{\mathrm{}}`$ is admissible if all its boundary circles are from vertex type pair of pants from the canonical decomposition and moreover, if one boundary circle from a vertex type pants is contained in $`\mathrm{\Sigma }`$ then the entire pants is contained in $`\mathrm{\Sigma }`$. In particular, $`𝒮_{\mathrm{}}\mathrm{\Sigma }`$ has no compact components.
2. Let $`\phi `$ be a homeomorphism of $`𝒮_{\mathrm{}}`$. One says that $`\phi `$ is asymptotically rigid if the following conditions are fulfilled:
* There exists an admissible subsurface $`\mathrm{\Sigma }_{g,n}𝒮_{\mathrm{}}`$ such that $`\phi (\mathrm{\Sigma }_{g,n})`$ is also admissible.
* The complement $`𝒮_{\mathrm{}}\mathrm{\Sigma }_{g,n}`$ is a union of $`n`$ infinite surfaces. Then the restriction $`\phi :𝒮_{\mathrm{}}\mathrm{\Sigma }_{g,n}𝒮_{\mathrm{}}\phi (\mathrm{\Sigma }_{g,n})`$ is rigid, meaning that it maps the pants decomposition into the pants decomposition and maps $`𝒯_{\mathrm{}}(𝒮_{\mathrm{}}\mathrm{\Sigma }_{g,n})`$ onto $`𝒯_{\mathrm{}}(𝒮_{\mathrm{}}\phi (\mathrm{\Sigma }_{g,n}))`$. Such a surface $`\mathrm{\Sigma }_{g,n}`$ is called a support for $`\phi `$.
One denotes by $`=(𝒮_{\mathrm{}})`$ the group of asymptotically rigid homeomorphisms of $`𝒮_{\mathrm{}}`$ up to isotopy and call it the asymptotic mapping class group of infinite genus.
In the same way one defined the asymptotic mapping class group $`(𝒮)`$, denoted by $``$ in .
###### Remark 1.1.
In genus zero (i.e. for the surface $`𝒮`$) a homeomorphism between two complements of admissible subsurfaces which maps the restrictions of the tree $`𝒯`$ one into the other is rigid, thus preserves the isotopy class of the pants decomposition. This is not anymore true in higher genus: the Dehn twist along a longitude preserves the edge-loop graph but it is not rigid, as a homeomorphism of the holed torus.
###### Remark 1.2.
Notice that, in general, rigid homeomorphisms $`\phi `$ do not have an invariant support i.e. an admissible $`\mathrm{\Sigma }_{g,n}`$ such that $`\phi (\mathrm{\Sigma }_{g,n})=\mathrm{\Sigma }_{g,n}`$. Take for instance a homeomorphism which translates the wrists along a geodesic ray in $`𝒯`$.
###### Remark 1.3.
Any admissible subsurface $`\mathrm{\Sigma }_{g,n}𝒮_{\mathrm{}}`$ has $`n=g+3`$. Moreover $`𝒮_{\mathrm{}}`$ is the ascending union $`_{g=1}^{\mathrm{}}\mathrm{\Sigma }_{g,g+3}`$. Instead of the wrist $`\mathrm{\Sigma }_{1,1}`$ use a surface of higher genus $`\mathrm{\Sigma }_{g,1}`$ and the same definitions as above. The admissible subsurfaces will be $`\mathrm{\Sigma }_{kg,k+3}`$. The asymptotic mapping class group obtained this way is finitely generated by small changes in the proof below.
###### Remark 1.4.
The surface $`𝒮_{\mathrm{}}`$ contains infinitely many compact surfaces of type $`(g,n)`$ with at least one boundary component. For any such compact subsurface $`\mathrm{\Sigma }_{g,n}𝒮_{\mathrm{}}`$, there is an obvious injective morphism $`i_{}:P(\mathrm{\Sigma }_{g,n})P`$. However, the morphism $`i_{}:(\mathrm{\Sigma }_{g,n})`$ is not always defined. Indeed, it exists if and only if the $`n`$ connected components of $`𝒮_{\mathrm{}}\mathrm{\Sigma }_{g,n}`$ are homeomorphic to each other, by asymptotically rigid homeomorphisms.
In particular, for any admissible subsurface $`\mathrm{\Sigma }_{g,n}`$ (hence $`n=g+3`$), $`i_{}`$ extends to an injective morphism $`i_{}:(\mathrm{\Sigma }_{g,n})`$ defined by rigid extension of homeomorphisms of $`\mathrm{\Sigma }_{g,n}`$ to $`𝒮_{\mathrm{}}`$.
#### 1.2.2 The group $``$ and the Thompson groups
###### Definition 1.4.
1. Let $`𝒯`$ be the planar trivalent tree. A partial tree automorphism of $`𝒯`$ is an isomorphism of graphs $`\phi :𝒯\tau _1𝒯\tau _2`$, where $`\tau _1`$ and $`\tau _2`$ are two finite trivalent subtrees of $`𝒯`$ (each vertex except the leaves are 3-valent). A connected component of $`𝒯\tau _1`$ or $`𝒯\tau _2`$ is a branch, that is, a rooted planar binary tree whose vertices are 3-valent, except the root, which is 2-valent. Each vertex of a branch has two descendant edges, and given an orientation to the plane, one may distinguish between the left and the right descendant edges. A partial automorphism $`\phi :𝒯\tau _1𝒯\tau _2`$ is planar if it maps each branch of $`𝒯\tau _1`$ onto the corresponding branch of $`𝒯\tau _2`$ by respecting the left and right ordering of the edges.
2. Two planar partial automorphisms $`\phi :𝒯\tau _1𝒯\tau _2`$ and $`\phi ^{}:𝒯\tau _1^{}𝒯\tau _2^{}`$ are equivalent, which is denoted $`\phi \phi ^{}`$, if and only if there exists a third $`\phi ^{\prime \prime }:𝒯\tau _1^{\prime \prime }𝒯\tau _2^{\prime \prime }`$ such that $`\tau _1\tau _1^{}\tau _1^{\prime \prime }`$, $`\tau _2\tau _2^{}\tau _2^{\prime \prime }`$ and $`\phi _{𝒯\tau _1^{\prime \prime }}=\phi _{𝒯\tau _1^{\prime \prime }}^{}=\phi ^{\prime \prime }`$.
3. If $`\phi `$ and $`\phi ^{}`$ are planar partial automorphisms, one can find $`\phi _0\phi `$ and $`\phi _0^{}\phi ^{}`$ such that the source of $`\phi _0`$ and the target of $`\phi _0^{}`$ coincide. The product $`[\phi ][\phi ^{}]=[\phi _0\phi _0^{}]`$ is well defined, as is easy to check. The set of equivalence classes of such automorphisms endowed with the above internal law, is a group with neutral element the class of $`id_𝒯`$. This is the Thompson group $`V`$.
###### Remark 1.5.
We warn the reader that our definition of the group $`V`$ is different from the standard one (as given in ). Nevertheless, the present group $`V`$ is isomorphic to the group denoted by the same letter in .
We introduce Thompson’s group $`T`$, the subgroup of $`V`$ acting on the circle (see ), which will play a key role in the proofs.
###### Definition 1.5 (Ptolemy-Thompson’s group $`T`$).
Choose a vertex $`v_0`$ of $`𝒯`$. Each $`gV`$ may be represented by a planar partial automorphism $`\phi :𝒯\tau _1𝒯\tau _2`$ such that $`v_0`$ belongs to $`\tau _1\tau _2`$. Let $`D_1`$ (respectively $`D_2`$) be a disk containing $`\tau _1`$ (respectively $`\tau _2`$), whose boundary circle $`S_1`$ (respectively $`S_2`$) passes through the leaves of $`\tau _1`$ (respectively $`\tau _2`$), giving to them a cycling ordering. If $`\phi `$ preserves this cycling ordering, which amounts to saying that the bijection from the set of leaves of $`\tau _1`$ onto the set of leaves of $`\tau _2`$ can be extended to an orientation preserving homeomorphism from $`S_1`$ onto $`S_2`$, then any other $`\phi ^{}`$ equivalent to $`\phi `$ also does, and one says that $`g`$ itself is circular. The subset of circular elements of $`V`$ is a subgroup, called the Ptolemy-Thompson group $`T`$.
###### Proposition 1.1.
Set $`P`$ for the inductive limit of the pure mapping class groups of admissible subsurfaces of $`𝒮_{\mathrm{}}`$. We have then the following exact sequences:
$$1PV1.$$
###### Proof.
Let $`\phi `$ be an asymptotically rigid homeomorphism of $`𝒮_{\mathrm{}}`$ and $`\mathrm{\Sigma }_{g,n}`$ a support for $`\phi `$. Then it maps $`𝒯_{\mathrm{}}(𝒮_{\mathrm{}}\mathrm{\Sigma }_{g,n})`$ onto $`𝒯_{\mathrm{}}(𝒮_{\mathrm{}}\phi (\mathrm{\Sigma }_{g,n}))`$, hence $`𝒯(𝒮_{\mathrm{}}\mathrm{\Sigma }_{g,n})`$ onto $`𝒯(𝒮_{\mathrm{}}\phi (\mathrm{\Sigma }_{g,n}))`$ by forgetting the action on the edge-loops. This may be identified with a planar partial automorphism $`\varphi :𝒯\tau _1𝒯\tau _2`$. The map $`[\phi ][\varphi ]V`$ is a group epimorphism. The kernel is the subgroup of isotopy classes of homeomorphisms inducing the identity outside a support, and hence is the direct limit of the pure mapping class groups. ∎
###### Remark 1.6.
In we prove the existence of a similar short exact sequence relating $``$ to $`V`$, which splits over the Ptolemy-Thompson group $`T`$. It is worth noticing that the present extension of $`V`$ is not split over $`T`$.
## 2 The proof of theorem 1.1
### 2.1 Specific elements of $``$
Recall that $`𝒮_{\mathrm{}}`$ has a canonical pants decomposition, as shown in figure 1. We fix an admissible subsurface $`A=\mathrm{\Sigma }_{1,4}`$ which contains a central wrist and an admissible $`B=\mathrm{\Sigma }_{0,3}\mathrm{\Sigma }_{1,4}`$ which is not adjacent to the wrist.
Let us consider now the elements of $``$ described in the pictures below. Specifically:
* Let $`\gamma `$ be a circle contained inside $`B`$ and parallel to the boundary curve labeled 3. Let $`t`$ be the right Dehn twist around $`\gamma `$. This means that, given an outward orientation to the surface, $`t`$ maps an arc crossing $`\gamma `$ transversely to an arc which turns right as it approaches $`\gamma `$. The dashed arcs (also called seams) on the left hand side picture figure out the boundary of the visible side of $`B`$. Their images by $`t`$ are represented on the right hand side picture,
* $`\pi `$ is the braiding, acting as a braid in $`(\mathrm{\Sigma }_{0,3})`$, with the support $`B`$. It rotates the circles 1 and 2 in the horizontal plane (spanned by the circles) counterclockwise.
Assume that $`B`$ is identified with the complex domain $`\{|z|7,|z3|1,|z+3|1\}`$. A specific homeomorphism in the mapping class of $`\pi `$ is the composition of the counterclockwise rotation of $`180`$ degrees around the origin — which exchanges the small boundary circles labeled 1 and 2 in the figure — with a map which rotates of $`180`$ degrees in the clockwise direction each boundary circle. The latter can be constructed as follows.
Let $`A`$ be an annulus in the plane, which we suppose for simplicity to be $`A=\{1|z|2\}`$. The homeomorphism $`D_{A,C}`$ acts as the counterclockwise rotation of $`180`$ degrees on the boundary circle $`C`$ and keeps the other boundary component pointwise fixed:
$$D_{A,C}(z)=\{\begin{array}{cc}z\mathrm{exp}(\pi \sqrt{1}(2|z|)),\hfill & \text{ if }C=\{|z|=1\}\hfill \\ z\mathrm{exp}(\pi \sqrt{1}(|z|1)),\hfill & \text{ otherwise}\hfill \end{array}$$
The map we wanted is $`D_{A_0,C_0}^1D_{A_1,C_1}^1D_{A_2,C_2}^1`$, where $`A_0=\{6|z|7\}`$, $`C_0=\{|z|=7\}`$, $`A_1=\{1|z3|2\}`$, $`C_1=\{|z3|=1\}`$, $`A_2=\{1|z+3|2\}`$, and $`C_2=\{|z+3|=1\}`$.
One has pictured also the images of the seams.
* $`\beta `$ is the order 3 rotation in the vertical plane of the paper. It is the unique globally rigid mapping class which permutes counterclockwise and cyclically the three boundary circles of $`B`$. An invariant support for $`\beta `$ is $`B`$.
* $`\alpha `$ is a twisted rotation of order 4 in the vertical plane which moves cyclically the labels of the boundary circles counterclockwise. Its support is a 4-holed torus $`A=\mathrm{\Sigma }_{1,4}`$.
Let $`\mathrm{\Sigma }_{0,5}`$ be the 5-holed sphere consisting of the union of $`B`$ with the edge pants $`Q`$ near $`B`$ and the next vertex pants $`B^{}`$ adjacent to $`Q`$. There are four boundary circles which are vertex type and one boundary circle which bounds a wrist. We perform first a rotation in $`^3`$ which preserves globally the pants decomposition and visible side, permutes counterclockwise and cyclically the four vertex type boundary circles of $`\mathrm{\Sigma }_{0,5}`$ and rotates the edge type circle according to one fourth twist. This rotation changes the position of the wrist $`\mathrm{\Sigma }_{1,1}`$ in $`^3`$. We consider next the clockwise rotation of this wrist alone, of angle $`\frac{\pi }{2}`$ around the vertical axis that meets the edge type circle in its center. This rotation restores the initial wrist position. The composition of the two partial rotations above is a homeomorphism of $`\mathrm{\Sigma }_{1,4}`$ that gives a well-defined element of $``$.
* Let $`a_1,b_1`$ be the meridian and longitude on the basic wrist in $`A`$. We denote by $`t_{a_1}`$ and $`t_{b_1}`$ the Dehn twists along these curves. Further, $`t_0`$ states for the Dehn twist $`t_0`$ along the boundary circle of the wrist.
###### Remark 2.1.
It is worthy to note that we have three types of Dehn twists: those along separating curves (conjugate either to $`t`$ (the boundary on the vertex pants) or with $`t_0`$ (the edge type pants) and those along non-separating curves which are conjugate to the twist around such a curve on the wrist.
### 2.2 Generators for $`P`$
Consider the following collection of simple curves drawn on $`𝒮_{\mathrm{}}`$:
Their description follows.
1. Choose, for each wrist, a longitude $`b_i`$, which turns once along the wrist.
2. For each pair of wrists we choose a circle joining them as follows. For each wrist we have an arc going from the base point of its attaching circle to the longitude and back to the opposite point of the circle. Then join these two pairs of points by a pair of parallel arcs in the horizontal surface, asking that the arc which joins the two base points be a geodesic path in the tree $`𝒯`$. We call them wrist-connecting loops.
3. Further we associate a loop to each pair consisting of a wrist and a vertex of the tree $`𝒯`$. A vertex gives rise to a pair of pants in $`𝒮_{\mathrm{}}`$. Two of the boundary components of these pants correspond to the directions to move away from the wrist. Thus we can define again an arc on the pair of pants which joins a point $`p`$ of the third circle (closest to the wrist), on the visible side of $`𝒮_{\mathrm{}}`$, to its opposite, on the hidden side, and separates the remaining two circles. Consider the loop resulting from gathering the following three kinds of arcs:
1. the arc on the wrist;
2. the arc on the pair of pants;
3. and a pair of parallel arcs which join them, asking that the arc which joins the point $`p`$ to the base point of the wrist be a geodesic path of the tree $`𝒯`$.
We call them vertex connecting loops.
4. Consider the loops that come from the canonical pants decomposition of $`𝒮`$ by doubling them. We call them the horizontal pants decomposition loops.
###### Lemma 2.1.
The set $``$ of Dehn twists along the meridians, the longitudes, the wrist connecting loops associated to edges, the vertex connecting loops and the horizontal pants decomposition loops generates $`P`$.
###### Proof.
It suffices to consider the finite case of an admissible surface with boundary that contains $`g`$ wrists and has $`g+3`$ boundary components. Then the lemma follows from , in which it is proved that the pure mapping class group of such a surface is generated by a set of Dehn twists $`_{g,n}`$ (with $`n=g+3`$). It suffices to check that all the Dehn twists belonging to $`_{g,n}`$ also belong to the set $``$. Referring to the notations of , there are four types of Dehn twists in $`_{g,n}`$: the $`\alpha _i`$’s, the $`\beta _i`$’s, the $`\gamma _{ij}`$’s and the $`\delta _i^{}s`$. The $`\alpha _i`$’s are associated to wrist-connecting or vertex connecting loops, the $`\beta _i`$’s are associated to longitudes $`b_i`$’s, the $`\gamma _{ij}`$’s are associated to wrist-connecting loops, except $`\gamma _{12}`$ which is associated to a vertex-connecting loop, and finally, the $`\delta _i`$’s are associated to the circles of the boundary of the surface, hence of the pants decomposition (after doubling them) of $`𝒮_{\mathrm{}}`$. Therefore all of them belong to $``$.
Remark also that in (, figure 1) the 1-handles are cyclically ordered and arranged on one side and then followed by all boundary components of the surface. However, we can arbitrarily permute the position of holes and 1-handles in the picture and keep the same system of generators. ∎
### 2.3 The action of $`𝕋`$ on the generators of $`P`$
#### 2.3.1 The groups $`𝕋`$ and $`T^{}`$
Consider the subgroup $`𝕋`$ of $``$ generated by the elements $`\alpha `$ and $`\beta `$. We will prove that the set of conjugacy classes for the action of $`𝕋`$ on $``$ is finite by considering the action of $`𝕋`$ on some planar subsurface of $`𝒮_{\mathrm{}}`$.
The surface obtained by puncturing (respectively deleting disjoint small open disks from) $`𝒮`$ at the midpoints of the edges is denoted by $`𝒮^{}`$ (and respectively $`𝒮^{}`$). The 2-dimensional thickening in $`𝒮`$ of the embedded tree $`𝒯`$ is an infinite planar surface, which will be called the visible side of $`𝒮`$, and will be denoted $`D`$. The intersection of $`D`$ with $`𝒮^{}`$ and $`𝒮^{}`$ is denoted $`D^{}`$ and $`D^{}`$, respectively.
The elements $`\alpha `$ and $`\beta `$ as defined above (i.e. as specific homeomorphisms, not only as mapping classes) keep invariant both $`𝒮^{}`$ and $`D^{}`$. If we crush the boundary circles to points then we obtain elements of $`(D^{})`$, and there is a well defined homomorphism $`𝕋(D^{})`$.
We studied in the asymptotic mapping class group $`(D^{})`$ denoted by $`T^{}`$ there. Recall from that:
###### Proposition 2.1.
The group $`T^{}`$ is generated by $`\alpha `$ and $`\beta `$.
This implies that $`𝕋T^{}`$ is an epimorphism. The relation between the asymptotic mapping class groups $`𝕋`$ and $`T^{}`$ is made precise by the following:
###### Lemma 2.2.
We have an exact sequence
$$0^{\mathrm{}}𝕋T^{}1$$
where the central factor $`^{\mathrm{}}`$ is the group of Dehn twists along attaching circles, normally generated by $`t_0`$.
###### Proof.
If an asymptotically rigid homeomorphism of $`𝒮_{\mathrm{}}`$ preserving $`D^{}`$ is isotopically trivial once the circles are crushed to points, then it is isotopic to a finite product of Dehn twists along those circles. Therefore, the kernel of $`𝕋T^{}`$ is contained in the subgroup denoted $`^{\mathrm{}}`$. Observe that $`\alpha ^4=t_0`$ in $`𝕋`$, so that $`t_0`$ belongs to the kernel of $`𝕋T^{}`$. Consequently, the kernel contains all the $`𝕋`$-conjugates of $`t_0`$, hence $`^{\mathrm{}}`$. ∎
Thus if we understand the action of $`T^{}`$ on the isotopy classes of arcs embedded in $`D^{}`$ then we can easily recover the action of $`𝕋`$ on homotopy classes of loops of $`𝒮^{}`$, up to some twists along attaching circles.
#### 2.3.2 The action of $`T^{}`$ on the isotopy classes of arcs of $`D^{}`$
The planar model of $`D^{}`$ is the punctured thick tree obtained from the binary tree by thickening in the plane and puncturing along midpoints of edges. The traces on $`D^{}`$ of the loops coming from the pants decomposition of $`𝒮`$ are arcs transversal to the edges. Thus $`D^{}`$ has a canonical decomposition into punctured hexagons. Each hexagon has three punctured sides coming from the arcs above, that we call separating side arcs. Moreover there are also three sides which are part of the boundary of $`D^{}`$ that we will call bounding side arcs. Notice that hexagons correspond to vertices of the binary tree, while separating side arcs. Further $`\beta `$ is the rotation of order 3 supported on the hexagon $`\overline{B}`$ (image of the pants $`B`$) and $`\alpha `$ is the rotation of order 4 that is supported on the union of $`\overline{B}`$ with an adjacent hexagon.
###### Lemma 2.3.
Let $`\gamma `$ be an arc embedded in $`D^{}`$ that joins two punctures. Then there exists some element of $`T^{}`$ that sends $`\gamma `$ in a prescribed arc joining the punctures 0 and 1.
###### Proof.
Recall from that the infinite braid group associated to the punctures $`B_{\mathrm{}}`$ is contained in $`T^{}`$. Further, there exists always a braid mapping class (supported in a compact subsurface of $`D^{}`$) sending the arc $`\gamma `$ in the prescribed one. ∎
###### Lemma 2.4.
The group $`T^{}`$ acts transitively on the set of separating side arcs.
###### Proof.
The group $`T^{}`$ contains $`PSL_2()`$, the group of orientation-preserving automorphisms of the tree $`𝒯`$, generated by $`\alpha ^2`$ and $`\beta `$. It acts transitively on the set of edges of $`𝒯`$, hence on the set of separating sides of the hexagons of $`D^{}`$.
An arc joining a puncture belonging to a hexagon $`H`$ to a bounding side of $`H`$ is called standard if it is entirely contained in $`H`$.
###### Lemma 2.5.
For any arc joining a puncture to a bounding side arc of a hexagon, there exists some element of $`T^{}`$ sending it into a standard arc joining the puncture $`0`$ to one of the bounding side of its hexagon.
###### Proof.
As above, $`PSL(2,)T^{}`$ also acts transitively on the set of all bounding sides. Thus we can use an element of $`T^{}`$ to send one end of our arc on a bounding side of the hexagon $`\overline{B}`$. Next, one composes by a braid element in $`B_{\mathrm{}}`$ that moves the other endpoint of the arc onto the puncture 0 and then makes the arc isotopic to a standard arc. ∎
Let $`t_{a_1}`$ and $`t_{b_1}`$ denote the Dehn twists along a meridian $`a_1`$ and a longitude $`b_1`$ on the wrist.
###### Corollary 2.1.
The elements $`\alpha ,\beta ,t,t_{a_1},t_{b_1},\pi `$, a Dehn twist along one wrist connecting loop and a Dehn twist along a vertex connecting loop generate $``$.
## 3 The rational homology of $``$
###### Theorem 3.1.
The rational homology of $``$ is isomorphic to the stable rational homology of the mapping class group: $`H_{}(,)H_{}(P,)`$.
###### Proof.
Recall first the theorem of stability, due to J. Harer (see ): Let $`R`$ be a connected subsurface of genus $`g_R`$ of a connected compact surface $`S`$ with at least one boundary component. Then the map $`H_n(P_R,)H_n(P_S,)`$ induced by the natural morphism $`P_RP_S`$ is an isomorphism if $`g_R2n+1`$.
The pure mapping class group $`P`$ is the inductive limit of the pure mapping class groups $`P_R`$, for all the compact subsurfaces $`R𝒮_{\mathrm{}}`$. It follows that $`H_n(P,)=\underset{\stackrel{}{R}}{lim}H_n(P_R,)=`$ $`H_n(P_R,)`$ for any compact subsurface $`R𝒮_{\mathrm{}}`$ of genus $`g_R2n+1`$. Therefore, the homology of $`P`$ is what is called the stable homology of the mapping class group. By Mumford’s conjecture proved in , $`H^{}(P,)`$ is isomorphic to $`[\kappa _1,\mathrm{},\kappa _i,\mathrm{}]`$, where $`\kappa _i`$, the $`i^{th}`$ Miller-Morita-Mumford class, has degree $`2i`$. Since $`H^{}(P,)=\mathrm{Hom}(H_{}(P,),)`$, each $`H_n(P,)`$ is finite dimensional over $``$.
Write now the Lyndon-Hochschild-Serre spectral sequence in homology associated with
$$1PV1$$
The second term is $`E_{p,q}^2=H_p(V,H_q(P,))`$. If we prove that $`V`$ acts trivially on the finite dimensional $``$-vector space $`H_q(P,)`$, and invoke a theorem of K. Brown () saying that $`V`$ is rationally acyclic, then the only possibly non-trivial term of the spectral sequence is $`E_{0,n}^2=H_n(P,)`$, and the proof is done.
Thus it remains to justify that $`V`$ acts trivially on the homology groups $`H_q(P,)`$, for any integer $`q0`$. This results from the fact that $`V`$ is not linear, as we explain below. Indeed, if $`dim_{}H_q(P,)=N`$, then $`\mathrm{Aut}(H_q(P,))GL(N,)`$. So, let $`\rho :VGL(N,)`$ be the representation resulting from the action of $`V`$ on $`H_q(P,)`$. Since $`V`$ is a simple group, $`\rho `$ is either trivial or injective. Suppose it is injective, so that $`V`$ is isomorphic to a finitely generated subgroup of $`SL(N,)`$. Now each finitely generated subgroup of $`SL(N,K)`$ for any field $`K`$ is residually finite. But $`V`$ is not residually finite, since its unique normal subgroup of finite index is the trivial subgroup. Therefore, $`\rho `$ is trivial.
###### Proposition 3.1.
The free universal mapping class group $``$ is perfect, and $`H_2(,)=`$. The generator of $`H^2(,)`$ is called the first universal Chern class of $``$, and is denoted $`c_1()`$.
###### Proof.
Recall that the pure mapping class group of a surface of type $`(g,n)`$ is perfect if $`g3`$. Consequently, $`P`$ is perfect. Since $`V`$ is perfect, $``$ is perfect as well.
The above spectral sequence may be written with integral coefficients. One obtains $`E_{2,0}^2=H_2(V,)=0`$ (see ), $`E_{1,1}^2=H_1(V,H_1(P,))=0`$ since $`P`$ is perfect, and $`E_{0,2}^2=H_0(V,H_2(P,))`$. By Harer’s theorem () and stability (), $`H_2(P,)`$. The action of $`V`$ on $`=H_2(P,)`$ must be trivial, since $`V`$ is simple, and it follows that $`E_{0,2}^2=`$. Thus, the only non-trivial $`E^{\mathrm{}}`$ term is $`E_{0,2}^{\mathrm{}}=E_{0,2}^2=`$, and this implies $`H_2(,)=`$. ∎
## 4 The symplectic representation in infinite genus
### 4.1 Hilbert spaces and symplectic structure associated to $`𝒮_{\mathrm{}}`$
There is a natural intersection form $`\omega :H_1(𝒮_{\mathrm{}},)\times H_1(𝒮_{\mathrm{}},)`$ on the homology of the infinite surface, but this is degenerate because it is obtained as a limit of intersection forms on surfaces with boundary. The $``$-module $`H_1(𝒮_{\mathrm{}},)`$ is the direct sum of two submodules: $`H_1(𝒮_{\mathrm{}},)=H_1(𝒮_{\mathrm{}},)_sH_1(𝒮_{\mathrm{}},)_{ns}`$, where $`H_1(𝒮_{\mathrm{}},)_s`$ is generated by the homology classes of separating circles of $`𝒮_{\mathrm{}}`$, while $`H_1(𝒮_{\mathrm{}},)_{ns}`$ is generated by the homology classes of non-separating circles of $`𝒮_{\mathrm{}}.`$ The kernel $`\mathrm{ker}\omega `$ of $`\omega `$ is $`H_1(𝒮_{\mathrm{}},)_s`$, and the restriction of $`\omega `$ to $`H_1(𝒮_{\mathrm{}},)_{ns}`$ is a symplectic form.
For each wrist torus occurring in the construction of $`𝒮_{\mathrm{}}`$ (see Definition 1.1), we consider the meridian $`a_k`$ and the longitude $`b_k`$, with intersection number $`\omega (a_k,b_k)=1`$. Note that both collections $`\{a_k,k\}`$ and $`\{b_k,k\}`$ are invariant by the mapping class group $`𝕋`$, since the generators $`\alpha `$ and $`\beta `$ rigidly map a wrist onto a wrist. Moreover, these collections are almost invariant by $``$, meaning that for each $`g`$, $`g(\{a_k,k\})`$ (respectively $`g(\{b_k,k\})`$) coincides up to isotopy with $`\{a_k,k\}`$ (respectively $`\{b_k,k\}`$) for all but finitely many elements.
The classes $`\{a_k,b_k,k\}`$ form a symplectic basis for $`H_1(𝒮_{\mathrm{}},)_{ns}`$. Each element of $``$ acts on $`H_1(𝒮_{\mathrm{}},)_{ns}`$ by preserving the intersection form $`\omega `$. In particular, there is a representation
$$\rho :P\mathrm{Sp}(2\mathrm{},)$$
where $`\mathrm{Sp}(2\mathrm{},)`$ is the inductive limit of the symplectic groups $`\mathrm{Sp}(2k,)`$, with respect to the natural inclusions $`\mathrm{Sp}(2k,)\mathrm{Sp}(2(k+1),)`$. Note, though, that if $`g`$ is not in $`P`$, it is not represented into $`\mathrm{Sp}(2\mathrm{},)`$, but into a larger symplectic group, that we are defining below.
One completes $`H_1(𝒮_{\mathrm{}},)_{ns}`$ as a real Hilbert space for which this basis orthonormal. Let $`_r`$ be this Hilbert space, and $`(.,.)`$ denote its scalar product.
Let $`J`$ be the almost-complex structure induced by $`\omega `$, i.e. the linear operator defined by $`\omega (v,w)=(v,Jw)`$ for all $`v,w`$ in $`_r`$. We have $`J^2=\mathrm{𝟏}`$. Each linear operator on $`_r`$ decomposes into a $`J`$-linear part $`T_1`$ and a $`J`$-antilinear part $`T_2`$, $`T=T_1+T_2`$, where $`T_1=\frac{TJTJ}{2}`$ and $`T_2=\frac{T+JTJ}{2}`$.
Recall that () the restricted symplectic group $`\mathrm{Sp}_{\mathrm{res}}(_r)`$ is defined as the group of symplectic (i.e. $`\omega `$ preserving) bounded invertible operators $`T`$ whose $`J`$-antilinear part $`T_2`$ is a Hilbert-Schmidt operator. An operator $`T`$ is called Hilbert-Schmidt if $`T_{HS}^2:=_iT(e_i)^2`$ is finite, where $`(e_i)_i`$ is an orthonormal Hilbert basis.
###### Theorem 4.1.
The symplectic representation of the mapping class group $`P`$ extends to a representation $`\widehat{\rho }:\mathrm{Sp}_{\mathrm{res}}(_r)`$ of $``$ into the restricted symplectic group.
### 4.2 Proof of Theorem 4.1
Instead of a direct proof we will introduce the complexification of $`_r`$ to be used also in the next section. Let $`=_r_{}`$. Extend $`\omega `$ and $`J`$ by $``$-linearity, and $`(.,.)`$ by sesquilinearity, and denote by $`\omega _{},J_{}`$ and $`(.,.)_{}`$ the extensions. Thus, $`(,(..)_{})`$ is a complex Hilbert space. Let $`B`$ be the indefinite hermitian form $`B(v,w)=\frac{1}{\sqrt{1}}\omega _{}(v,\overline{w})`$, for all $`v,w`$ in $``$, where $`\overline{w}`$ is the complex-conjugate of $`w`$.
Let $`\mathrm{Aut}(,\omega _{},B)`$ be the group of bounded invertible operators of $``$ which preserve $`\omega _{}`$ and $`B`$. The morphism $`\varphi :\mathrm{Sp}(_r)\mathrm{Aut}(,\omega _{},B)`$, given by $`\varphi (T)=T\mathrm{𝟏}_{}`$ is an isomorphism (see ), since any $`T\mathrm{Aut}(,\omega _{},)`$ commutes with the complex conjugation and hence stabilizes $`_r`$.
Since $`J_{}^2=id`$, $`=_+_{}`$ where $`_\pm =\mathrm{ker}(J\pm \sqrt{1}\mathrm{𝟏})`$. Moreover, the direct sum is orthogonal. The complex conjugation interchanges $`_+`$ and $`_{}`$. Let $`(e_k)_k`$ be an orthonormal basis of $`_+`$ and $`(f_k=\overline{e_k})_k`$ the conjugate basis of $`_{}`$.
According to (, 6.2) a symplectic operator $`T`$ belongs to $`\mathrm{Sp}_{\mathrm{res}}(_r)`$ if and only if the decomposition of $`\varphi (T)`$ relative to the direct sum $`_+_{}`$ in the basis $`(e_k)_k(f_k=\overline{e_k})_k`$ reads $`\left(\begin{array}{cc}\mathrm{\Phi }& \mathrm{\Psi }\\ \overline{\mathrm{\Psi }}& \overline{\mathrm{\Phi }}\end{array}\right)`$, where
1. $`{}_{}{}^{t}\overline{\mathrm{\Phi }}\mathrm{\Phi }{}_{}{}^{t}\mathrm{\Psi }\overline{\mathrm{\Psi }}=1`$ and $`{}_{}{}^{t}\overline{\mathrm{\Phi }}\mathrm{\Psi }={}_{}{}^{t}\mathrm{\Psi }\overline{\mathrm{\Phi }}`$, where $`{}_{}{}^{t}T`$ denotes the adjoint of $`T`$ with respect to $`(.,.)_{}`$.
2. $`\mathrm{\Psi }:_{}_+`$ is a Hilbert-Schmidt operator.
We will apply this criterion for the action of $``$.
Set
$$e_k=\frac{1}{\sqrt{2}}(a_k\sqrt{1}b_k),f_k=\frac{1}{\sqrt{2}}(a_k+\sqrt{1}b_k),$$
Then $`(e_k)_k`$ is an orthonormal basis of $`_+`$ and $`(f_k)_k`$ is the conjugate orthonormal basis of $`_{}`$. Moreover, $`\omega _{}(e_k,e_l)=\omega _{}(f_k,f_l)=0`$, $`\omega _{}(e_k,f_l)=\sqrt{1}\delta _{kl}`$, and $`B(e_k,e_l)=B(f_k,f_l)=\delta _{kl}`$, $`B(e_k,f_l)=0`$ for all $`k,l`$.
Consider now the action $`\widehat{\rho }(g)`$ of $`g`$ on the $``$-invariant subspace $`_r`$. We must check that $`\mathrm{\Psi }(\varphi (\widehat{\rho }(g)))`$ is a Hilbert-Schmidt operator. In fact, it is a finite rank operator. Let $`\mathrm{\Sigma }_{h,n}`$ be an admissible surface for $`g`$, that is $`g`$ is the mapping class of a homeomorphism $`G`$ so that $`G:𝒮_{\mathrm{}}\mathrm{\Sigma }_{h,n}𝒮_{\mathrm{}}\phi (\mathrm{\Sigma }_{h,n})`$ is rigid. Any wrist torus $`T_k=\mathrm{\Sigma }_{1,1}`$ of $`𝒮_{\mathrm{}}\mathrm{\Sigma }_{g,n}`$ is rigidly mapped by $`G`$ onto another corresponding wrist torus $`T_{\sigma (k)}`$, for some infinite permutation $`\sigma `$. Therefore, for any such $`T_k`$, the associated matrices are such that
$$\varphi (\widehat{\rho }(g))(e_k)=e_{\sigma (k)},\varphi (\widehat{\rho }(g))(f_k)=f_{\sigma (k)}$$
In particular, for all but finitely many $`f_k`$ (i.e. excepting those corresponding to tori $`T_k\mathrm{\Sigma }_{h,n}`$) we have $`\varphi (\widehat{\rho }(g))(f_k)_{}`$. Now $`\mathrm{\Psi }(\varphi (\widehat{\rho }(g))):_{}_+`$ corresponds to the components of $`\varphi (\widehat{\rho }(g))(f_k)`$ in $`_+`$. This means that $`\mathrm{\Psi }(\varphi (\widehat{\rho }(g)))`$ has finite rank, and in particular, it is Hilbert-Schmidt. This proves that $`\widehat{\rho }(g)\mathrm{Sp}_{\mathrm{res}}(_r)`$, as claimed.
## 5 The universal first Chern class
### 5.1 The Pressley-Segal extension
Let $``$ be a polarized separable Hilbert space as above, that is, the orthogonal sum of two separable Hilbert spaces $`=_+_{}`$. The restricted linear group $`\mathrm{GL}_{\mathrm{res}}()`$ (see ) is the Banach-Lie group of operators in $`\mathrm{GL}()`$ whose block decomposition $`A=\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)`$ is such that $`b`$ and $`c`$ are Hilbert-Schmidt operators. Moreover, the invertibility of $`A`$ implies that $`a`$ is Fredholm in $`_+`$, and has an index $`\mathrm{ind}(a)`$. This gives a homomorphism $`\mathrm{ind}:\mathrm{GL}_{\mathrm{res}}()`$, that induces an isomorphism $`\pi _0(\mathrm{GL}_{\mathrm{res}}())`$. Denote by $`\mathrm{GL}_{\mathrm{res}}^0()`$ the connected component of the identity. Then $`\mathrm{GL}_{\mathrm{res}}^0()`$ is a perfect group (cf. , §5.4).
###### Proposition 5.1.
The restricted symplectic group $`\mathrm{Sp}_{\mathrm{res}}(_r)`$ embeds into the restricted linear group $`\mathrm{GL}_{\mathrm{res}}^0()`$. It is given the induced topology.
###### Proof.
Let $`\left(\begin{array}{cc}\mathrm{\Phi }& \mathrm{\Psi }\\ \overline{\mathrm{\Psi }}& \overline{\mathrm{\Phi }}\end{array}\right)`$ be in $`\mathrm{Sp}_{\mathrm{res}}(_r)`$. Since $`\mathrm{\Psi }`$ is a Hilbert-Schmidt operator, $`K={}_{}{}^{t}\mathrm{\Psi }\overline{\mathrm{\Psi }}`$ is trace-class, hence compact. Then $`{}_{}{}^{t}\overline{\mathrm{\Phi }}\mathrm{\Phi }=1+K1`$, hence $`{}_{}{}^{t}\overline{\mathrm{\Phi }}\mathrm{\Phi }1`$ is injective, and the Fredholm alternative implies it is invertible. In particular, $`\mathrm{\Phi }`$ itself is invertible, and has null index. ∎
Pressley-Segal’s extension of the restricted linear group. Let $`_1(_+)`$ denote the ideal of trace-class operators of $`_+`$. It is a Banach algebra for the norm $`b_1=Tr(\sqrt{b^{}b})`$, where $`Tr`$ denotes the trace form. We say that an invertible operator $`q`$ of $`_+`$ has a determinant if $`qid__+=Q`$ is trace-class. Its determinant is the complex number $`det(q)=_{i=0}^+\mathrm{}Tr(^iQ)`$, where $`^iQ`$ is the operator of the Hilbert space $`^i_+`$ induced by $`Q`$ (cf. ).
Denote by $`𝔗`$ the subgroup of $`GL(_+)`$ consisting of operators which have a determinant, and by $`𝔗_1`$ the kernel of the morphism $`\mathrm{det}:𝔗^{}`$.
Let $`𝔈`$ be the subgroup of $`\mathrm{GL}_{\mathrm{res}}()\times \mathrm{GL}(_+)`$ consisting of pairs $`(A,q)`$ such that $`aq`$ is trace-class. Then $`\mathrm{ind}(a)=\mathrm{ind}(q+(aq))=\mathrm{ind}(q)=0`$, so that $`A`$ belongs to $`\mathrm{GL}_{\mathrm{res}}^0()`$. There is a short exact sequence
$$1𝔗\stackrel{i}{}𝔈\stackrel{p}{}\mathrm{GL}_{\mathrm{res}}^0()1$$
called the Pressley-Segal extension. Here $`p(A,q)=A`$, and $`i(q)=(\mathrm{𝟏}_{},q)`$. It induces the central extension
$$1\frac{𝔗}{𝔗_1}^{}\frac{𝔈}{𝔗_1}\mathrm{GL}_{\mathrm{res}}^0()1$$
The corresponding cohomology class in $`H^2(\mathrm{GL}_{\mathrm{res}}^0(),^{})`$ is denoted by $`PS`$, and called the Pressley-Segal class of the restricted linear group.
The Pressley-Segal class and the universal first Chern class. For $`G`$ a Banach-Lie group, set $`Ext(G,^{})`$ for the set of equivalence classes of central extensions of $`G`$ by $`^{}`$, which are locally trivial fibrations (see ). Note that $`Ext(G,^{})`$ must not be confused with the group of continuous cohomology $`H_{cont}^2(G,^{})`$, since the latter only classifies the topologically split central extensions. One introduces two maps $`H^2(G,^{})\stackrel{\delta }{}Ext(G,^{})\stackrel{\tau }{}H_{top}^2(G,)`$. The map $`\delta `$ associates to an extension $`E`$ in $`Ext(G,^{})`$ its cohomology class in the Eilenberg-McLane cohomology of $`G`$. The map $`\tau `$ is the composition of $`Ext(G,^{})[G,B^{}=K(,2)]`$, which sends an extension to the homotopy class of its classifying map $`GB^{}`$, with the isomorphism $`[G,K(,2)]H_{top}^2(G,)`$.
Let us apply this formalism to $`G=\mathrm{GL}_{\mathrm{res}}^0()`$ and the central Pressley-Segal extension, viewed as an element $`𝒫𝒮Ext(\mathrm{GL}_{\mathrm{res}}^0(,^{}))`$. Then $`\delta (𝒫𝒮)=PS`$. The point is that $`\mathrm{GL}_{\mathrm{res}}^0()`$ is a homotopic model of the classifying space $`BU`$. In fact $`𝔈`$ is contractible (see , 6.6.2) and $`𝔗`$ is homotopically equivalent to $`U`$ (see ), hence the claim. It follows that the fibration $`𝒫𝒮`$ corresponds to the universal first Chern class, that is, $`\tau (𝒫𝒮)=c_1(BU)H^2(BU,)`$.
### 5.2 Cocycles on $`\mathrm{GL}_{\mathrm{res}}^0()`$, $`\mathrm{Sp}_{\mathrm{res}}(_r)`$ and $``$
###### Lemma 5.1.
The class $`\iota ^{}(PS)`$ in $`H^2(\mathrm{Sp}_{\mathrm{res}}(_r),^{})`$ is represented by the cocycle
$$C_1(g,g^{})=det(\mathrm{\Phi }(g)\mathrm{\Phi }(g^{})\mathrm{\Phi }(gg^{})^1)$$
###### Proof.
Let $`𝒱`$ be the open subset of $`\mathrm{GL}_{\mathrm{res}}^0()`$ consisting of operators $`A`$ such that $`a`$ is invertible. It is known that the central Pressley-Segal extensions splits over $`𝒱`$, since it has the section $`\sigma :𝒱𝔈,`$ $`A(A,a)`$. In particular, there is a local cocycle for $`PS`$ given by the formula (, 6.6.4): $`C(A,A^{})=det(1+aa^{}a^{\prime \prime 1})`$, for $`A,A^{}𝒱`$, where $`a^{\prime \prime }`$ is the first block of $`AA^{}`$. It suffices now to observe that $`\mathrm{Sp}_{\mathrm{res}}(_r)`$ embeds into $`𝒱`$. ∎
In order to prove Proposition 5.2 below, we need the following result that contrasts sharply with the finite dimensional case:
###### Lemma 5.2.
The restricted symplectic group $`\mathrm{Sp}_{\mathrm{res}}(_r)`$ is contractible.
###### Proof.
Denote by $`Z`$ the set of symmetric Hilbert-Schmidt operators $`_{}_+`$ with norm $`<1`$. Clearly, $`Z`$ is a contractible subspace of the Banach space of Hilbert-Schmidt operators. The group $`\mathrm{Sp}_{\mathrm{res}}(_r)`$ acts transitively and continuously (see p. 177) on $`Z`$ by means of
$$g(S)=(\mathrm{\Phi }(g)S+\mathrm{\Psi }(g))(\overline{\mathrm{\Psi }(g)}S+\overline{\mathrm{\Phi }(g)})^1Z,\mathrm{for}g\mathrm{Sp}_{\mathrm{res}}(_r),SZ$$
The stabilizer of $`S=0`$ is the group of matrices $`\left(\begin{array}{cc}\mathrm{\Phi }& 0\\ 0& \overline{\mathrm{\Phi }}\end{array}\right)`$ such that $`\mathrm{\Phi }`$ is unitary in $`_+`$. Thus, it is isomorphic to $`𝒰(_+)`$. By a result of Kuiper (), $`𝒰(_+)`$ is contractible. The claim is now a consequence of the contractibility of $`\mathrm{Sp}_{\mathrm{res}}(_r)/𝒰(_+)Z`$. ∎
###### Proposition 5.2.
For each integer $`n`$, there is a well-defined continuous cocycle $`C_n`$ defined on $`\mathrm{Sp}_{\mathrm{res}}(_r)`$, with values in $`^{}`$, such that
$$C_n(g,g^{})=det\left((\mathrm{\Phi }(g)\mathrm{\Phi }(g^{})\mathrm{\Phi }(gg^{})^1)^{\frac{1}{n}}\right)$$
Moreover, $`\frac{C_n}{|C_n|}`$ may be lifted to a real cocycle $`\widehat{\varsigma _n}:\mathrm{Sp}_{\mathrm{res}}(_r)\times \mathrm{Sp}_{\mathrm{res}}(_r)`$ such that
$$\frac{C_n(g,g^{})}{|C_n(g,g^{})|}=e^{2i\pi \widehat{\varsigma _n}(g,g^{})},\mathrm{for}\mathrm{all}g,g^{}\mathrm{Sp}_{\mathrm{res}}(_r)$$
The restriction $`\varsigma _1`$ of $`\widehat{\varsigma _1}`$ to $`\mathrm{Sp}(2\mathrm{},)`$ defines an integral cohomology class $`[\varsigma _1]H^2(\mathrm{Sp}(2\mathrm{},),)`$.
###### Proof.
In fact, $`\mathrm{\Phi }(g)^1\mathrm{\Phi }(g^{})^1\mathrm{\Phi }(gg^{})=1+(\mathrm{\Phi }(g^{})^1\mathrm{\Phi }(g)^1\mathrm{\Psi }(g)(\overline{\mathrm{\Psi }(g^{})})`$. But, according to (, p. 168) we have
$$\mathrm{\Phi }^1(g)\mathrm{\Psi }(g)<1\mathrm{and}\overline{\mathrm{\Psi }(g^{})}\mathrm{\Phi }(g^{})^1<1$$
Thus, there is a non-ambiguous definition of $`(\mathrm{\Phi }^1(g)\mathrm{\Phi }(gg^{})\mathrm{\Phi }^1(g^{}))^{\frac{1}{n}}`$ given by an absolutely convergent series.
The existence of $`\varsigma _n`$ is now an immediate consequence of the preceding lemma.
The map $`\mathrm{}:g\mathrm{Sp}(2\mathrm{},)\mathrm{}(g)=\frac{det(\mathrm{\Phi }(g))}{det(\mathrm{\Phi }(g))}`$ is well-defined, so that the cocycle
$$(g,g^{})\mathrm{Sp}(2\mathrm{},)\times \mathrm{Sp}(2\mathrm{},)e^{2i\pi \varsigma _1(g,g^{})}$$
is the coboundary of $`\mathrm{}`$. This proves that the cohomology class of $`\varsigma _1`$ restricted to $`\mathrm{Sp}(2\mathrm{},)`$ is integral. ∎
###### Remark 5.1.
1. The restrictions of the real cocycles $`\varsigma _n`$ on the finite dimensional Lie group $`\mathrm{Sp}(2g,)`$ are those constructed by Dupont-Guichardet-Wigner (see ). In fact, the authors of proved that the cohomology class of the restriction of $`\varsigma _1`$ to $`\mathrm{Sp}(2g,)`$ is integral, and is the image in $`H^2(\mathrm{Sp}(2g,),)`$ of the generator of $`H_{\mathrm{bor}}^2(\mathrm{Sp}(2g,),)=`$, the second group of borelian cohomology of $`\mathrm{Sp}(2g,)`$. They prove also that it is the image of the first Chern class $`c_1(BU(g,))`$ by the composition of maps
$$H^2(BU(g,),)H^2(B\mathrm{Sp}(2g,),)H^2(B\mathrm{Sp}(2g,)^\delta ,)H^2(\mathrm{Sp}(2g,),)H^2(\mathrm{Sp}(2g,),),$$
where $`B\mathrm{Sp}(2g,)^\delta `$ is the classifying space of $`\mathrm{Sp}(2g,)`$ as a discrete group.
2. The remark above implies that the map
$$H^{}(BU,)H^{}(B\mathrm{Sp}(2\mathrm{},),)H^{}(\mathrm{Sp}(2\mathrm{},),).$$
sends the first universal Chern class $`c_1(BU)`$ onto $`[\varsigma _1]H^2(\mathrm{Sp}(2\mathrm{},),).`$ Further, the symplectic representation $`\rho :P\mathrm{Sp}(2\mathrm{},)`$ maps $`[\varsigma _1]`$ onto the generator $`c_1(P)`$ of $`H^2(P,)`$.
3. According to (, Theorem 6.2.3), the Berezin-Segal-Shale-Weil cocycle is the complex conjugate of the cocycle $`C_{\frac{1}{2}}`$.
###### Theorem 5.1.
Let $`[\widehat{\varsigma _1}]H^2(\mathrm{Sp}_{\mathrm{res}}(_r),)`$ be the cohomology class of $`\widehat{\varsigma _1}`$. The pull-back of $`[\widehat{\varsigma _1}]`$ in $`H^2(,)`$ by the representation $`\widehat{\rho }`$ of Theorem 1.3 is integral, and is the natural image of the generator $`c_1()`$ of $`H^2(,)`$ in $`H^2(,)`$.
###### Proof.
Let $`\iota :P`$ and $`j:\mathrm{Sp}(2\mathrm{},)\mathrm{Sp}_{\mathrm{res}}(_r)`$ be the natural embeddings. Plainly, $`j\rho =\widehat{\rho }\iota `$. Since $`j^{}:H^2(\mathrm{Sp}_{\mathrm{res}}(_r),)H^2(\mathrm{Sp}(2\mathrm{},),)`$ maps $`[\widehat{\varsigma _1}]`$ onto $`[\varsigma _1]`$, one has $`\iota ^{}(\widehat{\rho }^{}[\widehat{\varsigma _1}])=\rho ^{}[\varsigma _1]`$. Let us denote by $`\overline{c_1}(P)`$ (respectively $`\overline{c_1}()`$) the image of $`c_1(P)`$ (respectively $`c_1()`$) in $`H^2(P,)`$ (respectively $`H^2(,)`$). According to Remark 5.1, 2., $`\rho ^{}[\varsigma _1]=\overline{c_1}(P)`$. By Proposition 3.1, $`\iota ^{}(\overline{c_1}(P))=\overline{c_1}()`$, hence $`\widehat{\rho }^{}[\widehat{\varsigma _1}]=\overline{c_1}()`$. ∎ |
warning/0506/math0506097.html | ar5iv | text | # The Parametric Degree of a Rational Surface
## Introduction
The parametric degree of a rational surface is the degree of the polynomials in the smallest possible proper parametrization. In the absence of base points, the parametric degree is just the square root of the degree of the surface. On the other hand, there are examples of series of rational surfaces showing that the number of base points in the smallest parametrization can be much larger than the degree (see Example 6). It is therefore clear that the degree does not tell too much about the parametric degree of a rational surface.
Example 1 below shows that the parametric degree is not a geometric but an arithmetic concept, in the sense that it depends on the choice of the ground field. For the complex case, gave upper and lower estimations for the parametric degree in terms of the degree of the surface.
In this paper, we introduce two geometrical invariants of a rational surface, namely level and keel. These two numbers govern the parametric degree in the sense that there exist linear upper and lower bounds. In particular, level and keel determine the parametric degree up to a multiplicative factor of $`2`$, independent on the choice of the ground field (as long as it is perfect). We can therefore say that the parametric degree depends “just slightly” on the choice of the ground field, as this choice can change by a factor of at most 2. Of course, this is a bit misleading because there are surfaces for which rationality depends on the choice of the field, for instance surfaces that are rational over $``$ but not rational over $``$.
We note that this result implies that the question of rationality is decidable over any perfect field with decidable first order theory: first, we compute a parametrization over the algebraic closure, using . If it is of degree $`d`$, then, by the result in this paper, the surface is rational iff there is a parametrization of degree at most $`2d`$. For a fixed surface, this question can be formulated as a first order sentence. Of course, decidability of rationality of surfaces is well-known for the real case by Comesatti’s theorem (see ).
## 1 Parametric Degree
Throughout the paper, we fix field $`𝕂`$, which is assumed to be perfect. A proper parametrization of a rational surface is a birational map $`\nu :^2S^r`$, $`r2`$, which is defined over $`𝕂`$. We can write $`\nu `$ as
$$(x_0:x_1:x_2)(F_0(x_0,x_1,x_2):\mathrm{}:F_r(x_0,x_1,x_2)),$$
with $`F_0,\mathrm{},F_r𝕂[x_0,x_1,x_2]`$ homogeneous of the same degree $`d`$ and without a common divisor. This representation is unique up to the multiplication by a nonzero constant. The number $`d`$ is uniquely defined and it is called the degree of the parametrization. (This should not be mixed up with the concept of degree of a rational map; in fact, any proper parametrization is a rational map of degree 1.)
The surface $`S`$ has, in general, proper parametrizations of different degree. The smallest possible degree is called the parametric degree of $`S`$, and denoted by $`pdeg(S)`$. For instance, we have $`pdeg(^2)=1`$, and $`pdeg(Q)=2`$ when $`Q`$ is a quadric surface in $`^2`$ with a $`𝕂`$-rational nonsingular point. In the second case, the inverse of the stereographic projection from this $`𝕂`$-rational point defines a parametrization of degree 2.
The parametric degree is an arithmetic concept, i.e. it depends on the choice of the field $`𝕂`$.
###### Example 1.
The torus with equation
$$(x^2+y^2+z^2+\frac{16}{25}w^2)^24x^2w^24y^2w^2=0$$
has a complex parametrization of degree 3, namely
$$(s:t:u)((s^2+u^2)(3u8t):i(s^2u^2)(3u8t):$$
$$6is(u^26tu+8t^2:20st(u3t)).$$
The smallest real parametrization has degree 4. So we have $`pdeg(S_{})=3`$ and $`pdeg(S_{})=4`$.
## 2 Preliminaries
Most definitions and theorems in this section are well-known. The single exception is the definition of the function $`\mathrm{nmc}`$ at the end of the section, which turns out to be convenient for the definition of the keel.
Let $`S`$ be a nonsingular projective surface over $`𝕂`$.
* A prime divisor is an irreducible curve on $`S`$. Note that it is required that the curve is defined over $`𝕂`$, but it may split into several components over the algebraic closure.
* $`\mathrm{Div}(S)`$, the group of divisors of $`S`$, is the free abelian group generated by the prime divisor. Divisors are denoted by captital letters.
* A divisor is called effective iff all its coefficients are nonnegative. If $`BA`$ is effective, then we also say that $`A`$ divides $`B`$ or $`AB`$.
* Two effective divisor without common component are equivalent iff they are two fibers of a rational map $`S^1`$. Linear equivalence is the finest equivalence relation on $`\mathrm{Div}(S)`$ which is compatible with addition and for which the previous statement is true.
* $`\mathrm{Pic}(S)`$ is the group of classes of divisors. Classes are denoted by capital letters.
* The class $`D`$ is called effective iff it has an effective divisor.
* $`|D|`$ is the set of all effective divisors in the class $`D`$. It has a natural structure of a projective space over $`𝕂`$. The dimension of this projective space is denoted by $`dim(D)`$. A linear system of divisors is a subset of $`|D|`$ corresponding to a projective subspace.
* If the point $`pS`$ is contained in all divisors in a linear system $`l`$, then $`p`$ is a base point of $`l`$. If the complete linear system $`|D|`$ has no base points, then we say that $`D`$ is free.
* If $`l`$ is a non-empty linear system, then the associated rational map is denoted by $`\varphi _l:S^{dim(l)}`$. It is defined outside the base locus. The codomain $`^{dim(l)}`$ is naturally identified with the dual projective space of $`l`$; the image of $`p`$ corresponds to the subset of divisors in $`l`$ passing through $`p`$, which is a hyperplane.
* The intersection product $`\mathrm{Pic}(S)^2`$ is symmetric and bi-additive, and if the classes $`A`$ resp. $`B`$ contain two effective divisors $`A_0`$ resps. $`B_0`$ without common component, then $`AB`$ is the number of common points of $`A_0`$ and $`B_0`$, properly counted. Especially, $`AB0`$ in this case.
* $`K`$ or $`K_S`$ is the canonical class of $`S`$.
* A class $`D`$ is called nef iff $`DC0`$ for all effective $`C`$.
We also need the following well-known theorems.
###### Theorem 1.
A birational regular map $`\varphi :S_1S_2`$ induces two homomorphisms, the pushforward $`\varphi _{}:\mathrm{Div}(S_1)\mathrm{Div}(S_2)`$ and the pullback $`\varphi ^{}:\mathrm{Div}(S_2)\mathrm{Div}(S_1)`$. Both functions are well defined on classes and preserve effectivity.
For $`C\mathrm{Pic}(S_1)`$ and $`D\mathrm{Pic}(S_2)`$, the following hold.
* $`(\varphi _{}\varphi ^{})(D)=D`$.
* $`dim(\varphi _{}(C))dim(C)`$.
* $`dim(\varphi ^{}(D))=dim(D)`$.
* $`\varphi _{}(C)D=C\varphi ^{}(D)`$.
* $`(\varphi _{}(C))^2C^2`$.
* $`(\varphi ^{}(D))^2=D^2`$.
* $`\varphi _{}(K_{S_1})=K_{S_2}`$.
###### Theorem 2.
Let $`S`$ be a nonsingular projective surface. Let $`E\mathrm{Div}(S)`$ be a prime divisor such that $`E^2=EK<0`$. A prime divisor with these properties is called exceptional divisor.
Then there exists a regular birational map $`\pi :SS^{}`$, called the blowing down of $`E`$, such that the kernel of $`\pi _{}`$ is generated by $`E`$. Moreover, $`K_S=\pi ^{}(K_S^{})+E`$.
Any birational regular map is a composition of such blowing down maps.
We define a function $`\mathrm{nmc}`$ (for “number of moving components”) from effective classes to nonnegative integers. Let $`X(D)^{dim(D)}`$ be the image of the associated rational map $`\varphi _D`$.
* If $`X(D)`$ is a point (this is the case iff $`dim(D)=0`$), then $`\mathrm{nmc}(D):=0`$.
* If $`X(D)`$ is a curve of degree $`m`$, then $`\mathrm{nmc}(D):=m`$.
* If $`X(D)`$ is a surface, then $`\mathrm{nmc}(D):=1`$.
## 3 Level and Keel
We first introduce level and keel for divisors on nonsingular surfaces. Then the concepts are transferred to embedded surfaces with arbitrary singularities, using a resolution of singularities.
Let $`D`$ be an effective divisor class of $`S`$. The level of $`D`$ is the supremum of all rational numbers $`p/q`$, $`q>0`$, such that $`qD+pK`$ is effective. If the supremum is assumed, then the keel of $`D`$ is equal to the supremum of all numbers of the form $`\frac{\mathrm{nmc}(qD+pK)}{q}`$ where $`p/q`$ is the level. If the supremum in the definition of the level is not assumed (for instance if the level is irrational or infinity), then the keel is defined as $`0`$.
###### Remark 1.
If some multiple of $`K`$ is effective, then we have $`\mathrm{level}(D)=\mathrm{}`$ and $`\mathrm{keel}(D)=0`$. Hence level and keel are only interesting if the Kodaira dimension of $`S`$ is negative.
###### Remark 2.
The numbers $`dim(D)`$ and $`\mathrm{nmc}(D)`$ are preserved under extension of the ground field. It follows that level and keel are preserved under extension of the ground field (i.e. they are geometric).
###### Example 2.
Let $`S=^2`$ and $`D=nL`$, where $`L`$ is the class of lines and $`n0`$. Then $`K=3L`$, and a class $`mL`$ is effective iff $`m0`$, and we have $`\mathrm{nmc}(0)=0`$. It follows that $`\mathrm{level}(D)=n/3`$ and $`\mathrm{keel}(D)=0`$.
###### Example 3.
Let $`S=^1\times ^1`$ and $`D=mF_1+nF_2`$, where $`F_1,F_2`$ are the classes of the fibers of the two projections and $`0mn`$. Then $`K=2F_12F_2`$, and a class $`aF_1+bF_2`$ is effective iff $`a0`$ and $`b0`$, and we have $`\mathrm{nmc}(aF_2)=a`$. It follows that $`\mathrm{level}(D)=m/2`$ and $`\mathrm{keel}(D)=nm`$.
###### Example 4.
We sketch an example which generalizes the two above.
Let $`\mathrm{\Gamma }`$ be a convex lattice polygon, i.e. the convex hull of a finite number of points in the plane with integer coordinates. The polygon $`\mathrm{\Gamma }`$ defines a nonsingular toric surface $`S`$ (the minimal resolution of the toric surface defined by the inner normals) and an effective divisor (the inverse image of the class of hyperplane section in the projective embedding defined by $`\mathrm{\Gamma }`$).
The class $`qD+pK`$ corresponds to the convex figure obtained by scaling $`\mathrm{\Gamma }`$ by a factor of $`q`$ and moving each edge $`p`$ steps inward. The class is effective iff this figure is non-empty. Hence the level is equal to $`p/q`$ if we can enlarge $`\mathrm{\Gamma }`$ by a factor of $`q`$, pass $`p`$ times to the convex hull of the interior points, and obtain a line segment or a point (see figure 1; cf also ). The keel is the number of points on this line segment or point, minus 1, divided by $`q`$.
We now define level and keel of a projective surface $`S^r`$ (possibly with singularities). Let $`\pi :\stackrel{~}{S}S`$ be a resolution of singularities, i.e. a proper birational map such that $`\stackrel{~}{S}`$ is nonsingular. Let $`H\mathrm{Pic}(\stackrel{~}{S})`$ be the class of the pullbacks of hyperplane sections. We define the level and the keel of $`S`$ as the level and the keel of $`H`$. Theorem 4 below says that this is independent of the choice of the desingularization.
###### Lemma 3.
Let $`\varphi :S_1S_2`$ be a regular birational map. Let $`D\mathrm{Pic}(S_2)`$. Let $`C`$ be an effective divisor of $`S_1`$ such that $`\varphi _{}(C)=0`$. Then $`C`$ is a common divisor of the linear system $`|\varphi ^{}(D)+C|`$.
###### Proof.
The proof proceeds by induction on the number of blowing downs into which the birational regular map $`\varphi `$ can be decomposed. First, assume that $`\varphi :S_1S_2`$ is the blowing down of an exceptional divisor $`E`$. Then $`C=nE`$ for some $`n0`$. Let $`m0`$ be the largest number such that $`mE`$ is a common divisor of $`|\varphi ^{}(D)+nE|`$. Then there exists an effective divisor in $`|\varphi ^{}(D)+(nm)E|`$ that does not have $`E`$ as component. It follows that $`(\varphi ^{}(D)+(nm)E)E=mn0`$, which shows that $`nE`$ is a common divisor.
Second, assume that $`\varphi `$ can be decomposed as $`S_1\stackrel{\varphi _1}{}S_3\stackrel{\varphi _2}{}S_2`$, where $`\varphi _1`$ is the blowing down of an exceptional divisor $`E`$. Because the pushforward preserves effectivity, $`\varphi _{1}^{}{}_{}{}^{}(C)`$ is effective. By induction, $`\varphi _{1}^{}{}_{}{}^{}(C)`$ is a common divisor of the linear system $`|\varphi _{2}^{}{}_{}{}^{}(D)+\varphi _{1}^{}{}_{}{}^{}(C)|`$. Because the dimension of a linear system is preserved by pullback, the equation $`dim(\varphi _{2}^{}{}_{}{}^{}(D)+\varphi _{1}^{}{}_{}{}^{}(C))=dim(\varphi _{2}^{}{}_{}{}^{}(D))`$ implies the equation $`dim(\varphi ^{}(D)+\varphi _{1}^{}{}_{}{}^{}\varphi _{1}^{}{}_{}{}^{}(C))=dim(\varphi ^{}(D))`$, which implies that $`\varphi _{1}^{}{}_{}{}^{}\varphi _{1}^{}{}_{}{}^{}(C)`$ is a common divisor of $`|\varphi ^{}(D)+\varphi _{1}^{}{}_{}{}^{}\varphi _{1}^{}{}_{}{}^{}(C)|`$. Now $`C\varphi _{1}^{}{}_{}{}^{}\varphi _{1}^{}{}_{}{}^{}(C)`$ lies in the kernel of $`\varphi _{1}^{}{}_{}{}^{}`$. Therefore it is a multiple $`nE`$ of the exceptional divisor.
We distinguish two cases. If $`n0`$, then $`C\varphi _{1}^{}{}_{}{}^{}\varphi _{1}^{}{}_{}{}^{}(C)`$, and it follows
$$dim(\varphi ^{}(D))dim(\varphi ^{}(D)+C)dim(\varphi ^{}(D)+\varphi _{1}^{}{}_{}{}^{}\varphi _{1}^{}{}_{}{}^{}(C))=dim(\varphi ^{}(D)),$$
hence we have equality everywhere and the statement is proved.
If $`n>0`$, then $`nE`$ is effective, and $`nE`$ is a common divisor of the linear system $`|\varphi _{1}^{}{}_{}{}^{}(\varphi _{2}^{}{}_{}{}^{}(D)+\varphi _{1}^{}{}_{}{}^{}(C))+nE|`$ by the induction base case. Therefore, we have
$$dim(\varphi ^{}(D)+C)=dim(\varphi _{1}^{}{}_{}{}^{}(\varphi _{2}^{}{}_{}{}^{}(D)+\varphi _{1}^{}{}_{}{}^{}(C))+nE)=$$
$$dim(\varphi _{1}^{}{}_{}{}^{}(\varphi _{2}^{}{}_{}{}^{}(D)+\varphi _{1}^{}{}_{}{}^{}(C)))=dim(\varphi ^{}(D)),$$
hence the statement is also proved. ∎
###### Theorem 4.
Let $`\pi _1:\stackrel{~}{S}_1S`$ and $`\pi _2:\stackrel{~}{S}_2S`$ be two desingularizations of $`S`$. Let $`H_1\mathrm{Pic}(\stackrel{~}{S}_1)`$ and $`H_2\mathrm{Pic}(\stackrel{~}{S}_2)`$ be the pullbacks of hyperplane sections. Then we have
$$\mathrm{level}(H_1)=\mathrm{level}(H_2),\mathrm{keel}(H_1)=\mathrm{keel}(H_2).$$
###### Proof.
First, let us assume that there exists a birational regular map $`\varphi :\stackrel{~}{S}_1\stackrel{~}{S}_2`$ such that $`\pi _1=\pi _2\varphi `$. Then $`\varphi `$ transforms hyperplane pullbacks to hyperplane pullbacks, i.e. $`\varphi ^{}(H_1)=H_2`$. Let $`p,q`$ be positive integers. By Theorem 2, the class $`C:=K_1\varphi ^{}(K_2)`$ is effective. By Lemma 3, the divisor $`pC`$ is a common divisor of the linear system $`|\varphi ^{}(qH_2+pK_2)+pC|=|qH_1+pK_1|`$. It follows that the two linear systems $`|qH_2+pK_2|`$ and $`|qH_1+pK_1|`$ have the same dimension and the same number of moving components, and the statement is proven.
In the general case, there exists a dominating desingularization $`\pi _3:\stackrel{~}{S}_3S`$ and birational regular maps $`\varphi _i:\stackrel{~}{S}_3\stackrel{~}{S}_i`$ such that $`\pi _3=\pi _i\varphi _i`$ for $`i=1,2`$. Hence it can be reduced to the special case above. ∎
###### Example 5.
Assume that $`\nu :^2S`$ is a parametrization of degree $`d`$ without base points. Then $`\nu `$ is regular and we can use it as resolution of singularities. It follows that $`H=dL`$ and we have $`\mathrm{level}(S)=\mathrm{level}(H)=d/3`$ and $`\mathrm{keel}(S)=\mathrm{keel}(H)=0`$, by Example 2.
It is also convenient to extend the notion of parametric degree to divisors. Let $`S`$ be a nonsingular surface. Let $`D\mathrm{Pic}(S)`$ be an nef divisor. A linear system $`l`$ of divisors is called parametrizing iff $`dim(l)=2`$ and $`\varphi _l:S^2`$ is birational. A class $`P`$ is called parametrizing iff $`|P|`$ contains a parametrizing linear system. Then we define $`pdeg(D)`$ as the minimum of all numbers $`PD`$, where $`P`$ is a parametrizing class.
###### Lemma 5.
Let $`S^r`$ be a (possibly singular) surface. Let $`\pi :\stackrel{~}{S}S`$ be a resolution of its singularities. Let $`H\mathrm{Pic}(\stackrel{~}{S})`$ be the class of the pullbacks of hyperplane sections. Then $`pdeg(S)=pdeg(H)`$.
###### Proof.
There is a one-to-one correspondence of parametrizations of $`S`$ and parametrizing linear systems of $`D`$, and the degrees coincide for corresponding parametrization/class. ∎
###### Remark 3.
Lemma 5 allows to reduce any relation between parametric degree, level and keel of a singular surfaces to the same relation between parametric degree, level and keel of a divisor on a nonsingular surface.
## 4 The Lower Bound
The main idea for establishing a lower bound for the parametric degree in terms of level and keel is to analyze what happens in the examples 2 and 5 when the parametrization has base points.
###### Theorem 6.
Let $`S^r`$ be a rational surface. Then we have
$$pdeg(S)3\mathrm{level}(S)+\mathrm{keel}(S).$$
###### Proof.
It suffices to prove the above inequality for a divisor $`H`$. We assume that $`PH=pdeg(H)`$ and that some linear subsystem of $`|P|`$ induces a birational map to $`^2`$. Moreover, we assume that $`l|P|`$ is a parametrizing linear system. If $`C`$ is a common divisor of $`l`$, then we would have $`(PC)HPH`$ because $`H`$ is nef. In this case, we can replace $`P`$ by $`PC`$ which is also a parametrizing class. As we can do this only finitely many times, because $`l`$ has only finitely many common components, we can assume that $`l`$ has no common components.
Claim 1: $`P`$ is nef. Indeed, if $`C`$ is a prime divisor, then $`CP0`$ because $`C`$ is not a common component of $`l`$.
Claim 2: if $`C`$ is a prime divisor with positive dimension, then $`PC>0`$. Indeed, we have $`P^2>0`$ because the image of $`\varphi _l`$ is $`^2`$, and therefore two generic divisors in $`l`$ intersect in a point outside the base locus. And $`|C|`$ contains two divisors without common component, hence $`C^20`$. Therefore $`PC`$ cannot be zero by the Hodge index theorem.
Claim 3: $`PK3`$. To prove this, we resolve the base points of $`\varphi _l:S^2`$ and get a birational regular map $`\pi :\stackrel{~}{S}S`$ such that $`\varphi _P\pi :\stackrel{~}{S}^2`$ is regular. Then $`\pi ^{}(P)`$ is nef by Theorem 2, and it follows
$$PK=\pi ^{}(P)\pi ^{}(K)=\pi ^{}(P)K_{\stackrel{~}{S}}\pi ^{}(P)(\varphi _l\pi )^{}(K_^2)$$
$$=(\varphi _l\pi )_{}(\pi ^{}(P))K_^2=L(3L)=3,$$
where $`L\mathrm{Pic}(^2)`$ is the class of lines.
For two positive integers $`p,q`$, the divisor $`qH+pK`$ can be effective only if $`(qH+pK)P=qpdeg(H)3p0`$. This proves that $`\mathrm{level}(H)pdeg(H)/3`$.
Now, assume $`p/q=\mathrm{level}(H)`$. We claim that $`\mathrm{nmc}(qH+pK)pdeg(H)q3p`$. Let $`F`$ be the greatest common divisor of $`qH+pK`$, and let $`B`$ be a generic divisor in $`|qH+pKF|`$. Then $`B`$ corresponds to a generic hyperplane section of the associated image $`X:=\varphi _{qH+pK}(S)`$. If the image $`X`$ is a point, then $`\mathrm{nmc}(qH+pK)=0`$ and the claim is true. If $`X`$ is a surface, then $`B`$ has positive dimension; it follows that $`pdeg(H)q3pBP1`$, and because $`\mathrm{nmc}(qH+pK)=1`$ the claim is true. If $`X`$ is a curve, then $`X`$ is necessarily rational, and $`\varphi _{qH+pK}`$ factors through a rational map $`S^1`$ which is associated to some divisor $`A`$. Because the divisors in $`A`$ are also the fibers of $`\varphi _{qH+pK}`$, we have $`B=mA`$, where $`m=\mathrm{nmc}(qH+pK)`$ is the number of intersection points of $`X`$ with a generic hyperplane. Since $`A`$ has positive dimension, we have
$$pdeg(H)q3pBP=mAPm=\mathrm{nmc}(qH+pK),$$
hence the claim is true also in this last case. This shows that
$$\mathrm{keel}(H)pdeg(H)3p/q=pdeg(H)\mathrm{level}(H).$$
###### Remark 4.
By analyzing Example 3 in a similar fashion, one can show that if $`S`$ has a parametrization of bidegree $`(m,n)`$, $`mn`$, then $`m2\mathrm{level}(S)`$ and $`n2\mathrm{level}(S)+\mathrm{keel}(S)`$.
###### Example 6.
Here is an example that shows how to use the concepts of level and keel in order to construct surfaces that have a high parametric degree.
Let $`n5`$ be an odd integer. Let $`S^3`$ be the surface given by the equation
$$z^{2n+1}x^2w^{2n1}y^nw^{n+1}=0.$$
We can compute the level and keel by using the parametrization
$$(s:t:u)((s^2u^{n2}+t^n)^ns:(s^2u^{n2}+t^n)^2tu^{n^22n}$$
$$:(s^2u^{n2}+t^n)u^{n^2n+1}:u^{n^2+1}).$$
By resolving the base points of the parametrization, we get a resolution $`\stackrel{~}{S}`$ of the singularities of $`S`$. Explicit analysis of the base points (see ) shows that there is one base point of multiplicity $`n^22n`$, with $`\frac{n3}{2}`$ base points of multiplicity $`2n`$ and $`2n+3`$ base points of multiplicity $`n`$ in the infinitely near, and one simple base point with $`n^22n1`$ simple base points in the infinitely near.
For positive integers $`p,q`$, the linear system $`|qH+pK|`$ on $`\stackrel{~}{S}`$ corresponds to the linear space of forms of degree $`q(n^2+1)3p`$ that vanish with multiplicity $`qrp`$ at each point of multiplicity $`r`$. Such forms exist for $`2p(2n+1)q`$, hence $`\mathrm{level}(S)=n+\frac{1}{2}`$. If $`2p=(2n+1)q`$, then the corresponding linear space is the vectorspace of forms of degree $`\frac{2n^26n1}{2}q`$ vanishing with multiplicity $`\frac{2n^26n1}{2}q`$ at the $`(n^22n)`$-fold base point and with multiplicity $`\frac{2n1}{2}q`$ at the $`\frac{n3}{2}`$ base points of multiplicity $`2n`$ in the infinitely near. Hence
$$\mathrm{keel}(S)=\frac{2n^26n1}{2}\frac{n3}{2}\frac{2n1}{2}=\frac{2n^25n5}{4}.$$
By Theorem 6, the parametric degree is greater than or equal to $`\frac{2n^2+7n+1}{4}`$. Because there is a parametrization of degree $`n^2+1`$, we know that the parametric degree grows proportional to the square of the implicit degree.
## 5 The Adjoint Chain
Adjoints are a tool for constructing minimal models of a given surface or higher-dimensional varieties. Starting with a nef divisor class, we keep alternating to blow down orthogonal exceptional divisors and adding the canonical class, until nefness does not hold any more. The last surface with nef class in the process is a minimal model with special properties. This technique has been used in to construct parametrizations in the case $`𝕂`$ is algebraically closed. Similar constructions appear in various other contexts, see for a survey.
Let $`S`$ be a nonsingular surface, and let $`D\mathrm{Pic}(S)`$. Following , we say that $`S`$ is $`D`$-minimal iff $`S`$ has no exceptional divisor orthogonal to $`D`$. We say $`S`$ is minimal iff $`S`$ is $`0`$-minimal. The following theorem is well-known (see ).
###### Theorem 7.
Let $`S`$ be a nonsingular surface and let $`D\mathrm{Pic}(S)`$. Then there exists a birational regular map $`\mu :SS_0`$, such that $`D=(\mu ^{}\mu _{})(D)`$ and $`S_0`$ is $`\mu _{}(D)`$-minimal. We call this a $`D`$-minimalization.
In the rest of the paper, we fix a nonsingular rational surface $`S`$ and a nef class $`D\mathrm{Pic}(S)`$ such that $`D^2>0`$ (for instance, the class of pullback of hyperplane sections in a resolution). The adjoint chain $`𝒮`$ is a chain of surfaces and birational regular maps
$$S\stackrel{\mu _0}{}S_0\stackrel{\mu _1}{}S_1\stackrel{\mu _2}{}\mathrm{}$$
and divisor classes $`D_i\mathrm{Pic}(S_i)`$, which is constructed recursively in the following way. First, we let $`\mu _0:SS_0`$ be a $`D`$-minimalization of $`S`$, and we let $`D_0:=\mu _{0}^{}{}_{}{}^{}(D_0)`$. Now assume that we have already defined $`S_i`$ and $`D_i`$. Let $`K_i`$ be the canonical class of $`S_i`$. If $`D_i+K_i`$ is not effective, then the adjoint chain ends; we denote the index of the last surface with $`a`$. Otherwise, we let $`\mu _{i+1}:S_iS_{i+1}`$ be a $`(D_i+K_i)`$-minimalization of $`S_i`$, and we let $`D_{i+1}:=(\mu _{i+1})_{}(D_i+K_i)`$. If the adjoint chain is infinite, then we set $`a:=\mathrm{}`$ (but we will prove that $`a`$ is finite).
###### Lemma 8.
The classes $`D_i`$ above are effective and nef. If $`i<a`$, then $`D_{i}^{}{}_{}{}^{2}>0`$.
###### Proof.
If $`𝕂`$ is algebraically closed, then the proof is well-known (, Lemma A.2 and Lemma A.3). There is only one step in the proofs of that uses the assumption that $`𝕂`$ is algebraically closed: in this case, any prime divisor $`C`$ of dimension 0 with $`CK<0`$ is exceptional, and we have $`C^2=CK=1`$ (Lemma A.1 in ). It follows that if $`D`$ is nef and $`D+K`$ is effective but not nef, then there exists an exceptional divisor orthogonal to $`D`$.
Here is an adaption of the proof to the case of non-closed fields: assume that $`D`$ is nef and $`D+K`$ is effective but not nef. Let $`C`$ be a prime divisor such that $`(D+K)C<0`$. Then $`CK<0`$, and $`dim(C)=0`$ because $`C`$ must be fixed in $`|D+K|`$. By the lemma below, $`C`$ is exceptional and $`CD`$ is an integral multiple of $`CK`$. This is only possible if $`CD=0`$, hence $`C`$ is orthogonal to $`D`$, and the rest of the proof works as in the case where $`𝕂`$ is algebraically closed. ∎
###### Lemma 9.
Let $`C`$ be a prime divisor such that $`dim(C)=0`$ and $`CK<0`$. Then $`C`$ is exceptional, and $`CD`$ is an integral multiple of $`CK`$ for all $`D\mathrm{Pic}(D)`$.
###### Proof.
Let $`\overline{S}`$ be the surface obtained by base field extension to the algebraic closure $`\overline{𝕂}`$. There are natural injections $`\mathrm{Div}(S)\mathrm{Div}(\overline{S})`$ and $`\mathrm{Pic}(S)\mathrm{Pic}(\overline{S})`$. In general, $`C`$ need not be a prime divisor in $`\mathrm{Div}(\overline{S})`$, but it has only simple components $`C=_{i=1}^rC_i`$. Each $`C_i`$ has dimension 0. Moreover, $`C_iD=C_jD`$ for any $`i,jr`$ and $`D`$ coming from $`S`$, because $`C_i`$ and $`C_j`$ are conjugate under the action of the Galois group of the extension $`𝕂\overline{𝕂}`$. Especially, $`C_iK=C_jK`$. It follows that $`C_iK<0`$. By Lemma A.1 in , each $`C_i`$ is exceptional, and $`C_i^2=C_iK=1`$ for all $`i`$. Hence $`CK=r`$, and $`CD`$ is an integral multiple of $`r`$.
It remains to show that $`C^2=r`$. For $`ij`$, we have $`dim(C_i+C_j)C_iC_j`$ by the Riemann-Roch theorem. On the other hand, $`dim(C_i+C_j)=0`$ because $`dim(C)=0`$, hence $`C_iC_j=0`$. Hence $`C^2=_{i=1}^rC_i^2=r`$. ∎
###### Lemma 10.
Let $`p,q,i`$ be integers such that $`ia`$ and $`pqi`$. Then the linear systems $`|qD+pK|`$ and $`|qD_i+(pqi)K_i|`$ have the same dimension and number of moving components.
###### Proof.
For $`j=0,\mathrm{},i`$, let $`\varphi _j:SS_j`$ be the map $`\mu _j\mathrm{}\mu _0`$. Then $`\varphi _j^{}(D_j)=\varphi _{j1}^{}(D_{j1}+K_{j1})`$. The class $`E_j:=\varphi _j^{}K_j\varphi _{j1}^{}K_{j1}`$ is effective. Therefore, the class
$$C:=qD+pK\varphi _i^{}(qD_i+(pqi)K_i)=\underset{j=0}{\overset{i}{}}(pqj)E_j$$
is effective, too. By Lemma 3, the unique divisor in $`|C|`$ is fixed in $`|qD+pK|`$. Because the pullback preserves the dimension and the number of moving components, the lemma follows. ∎
###### Corollary 11.
We have $`\mathrm{level}(D_i)=\mathrm{level}(D)i`$ and $`\mathrm{keel}(D_i)=\mathrm{keel}(D)`$.
###### Lemma 12.
We have $`a\mathrm{level}(D)`$. In particular, $`a`$ is finite.
###### Proof.
By Corollary 11, it suffices to prove that $`\mathrm{level}(D_a)0`$. But this is clear since $`D_a`$ is effective. ∎
The following lemmas can be used to compute level and keel in terms of the adjoint chain.
###### Lemma 13.
Assume that $`D_a=0`$. Then $`\mathrm{level}(D)=a`$ and $`\mathrm{keel}(D)=0`$.
###### Proof.
This follows immediately from Corollary 11 and from $`\mathrm{level}(0)=0`$ and $`\mathrm{keel}(0)=0`$. ∎
###### Lemma 14.
Assume that $`D_{a}^{}{}_{}{}^{2}=0`$ and $`D_a0`$. Then $`\mathrm{level}(D)=a`$ and $`\mathrm{keel}(D)=\mathrm{nmc}(D_a)=:k>0`$, and there is a free divisor $`P\mathrm{Pic}(S_a)`$ such that $`D_a=kP`$ and $`PK_a=2`$ and $`P^2=0`$ and $`dim(P)=1`$.
###### Proof.
The proof of Lemma A.7 in generalizes without problems to non-closed fields. This shows the existence of $`P`$ with the desired properties.
If $`p,q>0`$, then $`qD_a+pK_a`$ cannot be effective because $`(qD_a+pK_a)P=2p<0`$. Hence $`\mathrm{level}(D_a)=0`$ and $`\mathrm{level}(D)=a`$ by Corollary 11. Moreover, $`\mathrm{nmc}(qD_a)=\mathrm{nmc}(qkP)=qk`$, hence $`\mathrm{keel}(D_a)=\mathrm{keel}(D)=k`$. ∎
###### Lemma 15.
Assume that $`D_{a}^{}{}_{}{}^{2}>0`$. Then one of the following cases holds.
a) $`\mathrm{level}(D)=a+1/3`$, $`\mathrm{keel}(D)=0`$, and $`3D_a+K_a=0`$.
b) $`\mathrm{level}(D)=a+2/3`$, $`\mathrm{keel}(D)=0`$, and $`3D_a+2K_a=0`$.
c) $`\mathrm{level}(D)=a+1/2`$, $`\mathrm{keel}(D)=0`$, and $`2D_a+K_a=0`$.
d) $`\mathrm{level}(D)=a+1/2`$, $`\mathrm{keel}(D)=\mathrm{nmc}(2D_a+K_a)/2>0`$, and $`(2D_a+K_a)^2=0`$.
In particular, level and keel are rational numbers with a denominator dividing 6.
###### Proof.
Lemma A.8 from – which is also true in the case $`𝕂`$ is non-closed – says that we can conclude from $`D_{a}^{}{}_{}{}^{2}>0`$ that either $`3D_a+K_a=0`$, or $`3D_a+2K_a=0`$, or $`(2D_a+K_a)^2=0`$. Using Corollary 11, we get $`\mathrm{level}(D)=\mathrm{level}(D_a)+a=a+i/3`$ in the $`i`$-th case, for $`i=1,2`$. In the third case, we apply Lemma 13 or Lemma 14 to $`D:=2D_a+K_a`$, and we get either (c) or (d), depending whether $`2D_a+K_a`$ is zero or not. ∎
###### Remark 5.
At this point, it is instructive to revisit Example 4 again. Starting from a convex lattice polygon, we pass to the convex hull of the interior points $`a`$ times. If we obtain a single point, then $`D_a=0`$ holds for the corresponding toric surface. If we obtain a line segment with $`k+1`$ lattice points, then $`D_a=kP`$ for some $`P`$ with $`P^2=0`$ and $`PK_a=2`$, as in Lemma 14. If we obtain a lattice polygon without interior lattice points, then one of the following four cases holds:
a) after scaling by 3, we get a polygon with one interior point;
b) after scaling by 3 and passing to the convex hull of interior points, we get a polygon with one interior point;
c) after scaling by 2, we get a polygon with one interior point;
d) after scaling by 2, we get a polygon with several interior points that are all on a line.
Of course, these are instances of the 4 cases (a), (b), (c), (d) in Lemma 15.
###### Remark 6.
The lemmas above remind on the Kawamata Rationality Theorem and the Kawamata-Shokurov Base Point Free Theorem (see ): if $`D`$ is ample, then the “nefness value” $`v`$ is a rational number, and some multiple of $`D+vK`$ is free. The associated contraction morphism is either a blowing down, or a map with conic fibers, or a constant map. Of course, the Kawamata Rationality Theorem and the Kawamata-Shokurov Base Point Free Theorem hold in a much more general context (arbitrary dimension, rationality need not be assumed).
###### Remark 7.
If $`S`$ is a rational surface with degree $`d`$ and sectional genus $`p_1`$, then we have the inequality $`a+dim(D_a)p_1+\left(\genfrac{}{}{0pt}{}{2p_1d1}{2}\right)`$, by Lemma 8 in . Using the classification of surfaces occuring in Lemma 15 (see Lemma A.8 in ), it is easy to check that $`\mathrm{keel}(D)dim(D_a)`$ in all cases. Together with the upper bound $`p_1\left(\genfrac{}{}{0pt}{}{d1}{2}\right)`$ for the sectional genus, we get the bound $`\mathrm{level}(D)+\mathrm{keel}(D)d^4/2`$. Together with the bound in Theorem 20 below, we obtain the bound $`pdeg(S)3\mathrm{deg}(S)^4`$.
## 6 The Upper Bound
In order to establish an upper bound for the parametric degree, one has to construct a parametrization (or, equivalently, a parametrizing divisor class). The idea is to construct a minimal model using adjoints, and then to use the well-known classification of such minimal surfaces, due to Manin and Iskovskih .
###### Theorem 16.
Let $`S`$ be a minimal rational surface such that $`K`$ is nef and $`K^2>0`$. Then one of the following cases holds.
a) $`S^2`$; in this case, $`K=3L`$, where $`L`$ is the class of lines.
b) $`S`$ is isomorphic to a quadric in $`^3`$ or to the blowup of a singular quadric cone in $`^3`$ at its vertex. If $`Q`$ is the class of conic plane sections, the $`K=2Q`$.
c) $`S`$ is isomorphic to a Del Pezzo surface of degree 5 in $`^5`$. Its Picard group is cyclic, generated by $`K`$. The class of hyperplane sections is $`K`$.
d) $`S`$ is isomorphic to a Del Pezzo surface of degree 6 in $`^6`$. Its Picard group is again cyclic, generated by $`K`$, and the class of hyperplane sections is $`K`$.
###### Proof.
Let $`d:=K^2`$. By the classification of Del Pezzo surfaces over algebraically closed field (see , Theorem 24.4), we have $`1d9`$.
If $`d=9`$, then $`S`$ is a Severi-Brauer surface. As $`S`$ also has a parametrization over $`𝕂`$, it is isomorphic to $`^2`$ and (a) holds.
If $`d=8`$, then $`S`$ is isomorphic to a ruled surface $`F_n`$ over the algebraic closure $`\overline{𝕂}`$, where $`0n2`$. The case $`n=1`$ is not possible, because in this case $`S`$ would not be minimal; in the two remaining cases we have $`K=2Q`$ for some divisor $`Q`$, whose associated image is a quadric in $`^3`$.
If $`d=7`$, then $`S`$ is not minimal.
If $`d=6`$, then (d) holds.
If $`d=5`$, then (c) holds.
If $`d=4`$, then $`S`$ cannot be both minimal and rational over $`𝕂`$, by Theorem 1 in .
If $`d=3`$, $`2`$, or $`1`$, then $`S`$ cannot be both minimal and rational over $`𝕂`$, by Theorem 5.7 in . ∎
###### Theorem 17.
Let $`S`$ be a nonsingular rational surface and let $`D`$ be a nef divisor such that $`\mathrm{keel}(S)=0`$ and $`D^2>0`$. Then
$$pdeg(D)6\mathrm{level}(D).$$
###### Proof.
By lemmas 13, 14, 15, we can reduce to the case $`\mathrm{level}(D)=a`$ and $`D_a=0`$ by replacing $`D`$ by $`2D`$ or $`3D`$. Then $`K_a`$ is nef and $`(K_a)^2>0`$ because $`K_a`$ is the direct image of $`D_{a1}`$. Then $`S_a`$ satisfies the assumptions in Theorem 16. For each of the cases (a), (b), (c), we construct below a parametrizing class $`P\mathrm{Pic}(S_a)`$, such that $`P(K_a)6`$. Let $`\varphi _a:SS_a`$ be the minimalization map. Then $`\varphi _a^{}(P)`$ is a parametrizing divisor for $`S`$, and
$$\varphi _a^{}(P)D=P\varphi _{a}^{}{}_{}{}^{}(D)=P(D_aaK_a)6a.$$
Case (a): we take $`P:=L`$. Then $`P(K_a)=36`$.
Case (b): we take $`P:=Q`$. This is a parametrizing class because we can choose a parametrizing system $`l`$ as the linear system of conic sections through a fixed nonsingular point $`p`$ defined over $`𝕂`$. Such a point exists because $`S`$ is rational. The associated map is the stereographic projection from $`p`$, which is birational to $`^2`$. In this case, we have $`P(K_a)=2Q^2=46`$.
Case (c): we take $`P:=K_a`$, the class of hyperplane sections. As parametrizing system, we choose the set of all sections with hyperplanes containing the tangent plane through a fixed point $`p`$ defined over $`𝕂`$. The associated map is the projection from the tangent plane, which reduces the dimension by 3 and the degree by 4, hence it is birational to $`^2`$. In this case, we have $`P(K_a)=P^2=56`$.
Case (d): again we take $`P:=K_a`$. As parametrizing system, we choose the set of all sections with hyperplanes containing the tangent plane through a fixed point $`p`$ and through another fixed point $`q`$ outside the tangent plane, where both $`p`$ and $`q`$ are defined over $`𝕂`$. The associated map can be decomposed into the projection from the tangent plane, which is birational onto a quadric in $`^3`$, followed by the stereographic projection from the image of $`q`$. In this case, we have $`P(K_a)=P^2=6`$. ∎
###### Theorem 18.
Let $`S`$ be a nonsingular rational surface. Let $`P\mathrm{Pic}(S)`$ be a free class such that $`P^2=0`$ and $`PK=2`$ and $`dim(P)=1`$. Assume that $`S`$ is $`P`$-minimal. Then one of the following cases holds.
a) $`S`$ is isomorphic to the ruled surface $`F_n`$, $`n0`$. There exists an effective class $`C`$ such that $`CP=1`$, $`C^2=n`$, and $`K=(n2)P2C`$. The classes $`C`$ and $`P`$ generate $`\mathrm{Pic}(S)`$.
b) $`S`$ is isomorphic to the blowup of a nonsingular quadric at a point of degree 2 (i.e. defined over a quadratic extension of $`𝕂`$). The Picard group is generated by $`P`$ and the exceptional class $`E`$, the class of plane sections is $`P+E`$, the canonical class is $`2PE`$, and $`PE=2`$.
c) $`S`$ is isomorphic to the blowup of $`^2`$ at a point of degree 4. The Picard group is generated by the exceptional class $`E`$ and the class of lines $`L`$, we have $`P=2LE`$, and the canonical class is $`3L+E`$.
###### Proof.
Because $`P`$ is free, the associated map $`S^1`$ is regular, and $`P`$ is the class of fibers. The genus of a generic fiber is $`\frac{P^2+PK}{2}+1=0`$, hence the associated map gives $`S`$ the structur of a conic fibration. Let $`d:=K^2`$. Over the algebraic closure $`\overline{𝕂}`$ is a blowup of a ruled surface $`F_n`$ for some $`n`$, hence $`d8`$.
If $`d=8`$, then $`S`$ is minimal over $`\overline{𝕂}`$, and (a) holds.
If $`d=7`$, then $`S`$ is not minimal by Theorem 4.1 in .
If $`d=6`$, then (b) holds by Theorem 4.1 in .
If $`d=5`$, then (c) holds by Theorem 4.1 in .
If $`d=4`$, then $`S`$ cannot be both minimal and rational over $`𝕂`$, by Theorem 2 in .
If $`d=3`$, then $`S`$ cannot be both minimal and rational over $`𝕂`$, by Corollary 2.6 in .
If $`d=2`$ or $`1`$, then $`S`$ cannot be both minimal and rational over $`𝕂`$, by Corollary 1.7 in .
If $`d0`$, then $`S`$ cannot be both minimal and rational over $`𝕂`$, by Theorem 1.6 in . ∎
###### Theorem 19.
Let $`S`$ be a nonsingular rational surface and let $`D`$ be a nef divisor such that $`D^2>0`$ and $`\mathrm{keel}(S)>0`$. Then
$$pdeg(D)4\mathrm{level}(D)+2\mathrm{keel}(D).$$
###### Proof.
By lemmas 13, 14, 15, and by replacing $`D`$ by $`2D`$, we can reduce to the case $`\mathrm{level}(D)=a`$ and $`D_a=kP`$ for $`k=\mathrm{keel}(D)`$ and $`P`$ as in Lemma 14. Then we construct a parametrizing class $`Q`$ for each of the cases that arise in Theorem 18. If $`\varphi _a:SS_a`$ be the minimalization map, then $`\varphi _a^{}(Q)`$ is a parametrizing class for $`S`$; and we will prove the required upper bound for $`\varphi _a^{}(Q)D`$ in each case.
In case (a), we distinguish two subcases. If $`SF_0`$, then we take $`Q:=C+P`$. The image of the associated map is a ruled quadric in $`^3`$. Because a stereographic projection from a point defined over $`𝕂`$ is birational onto the plane, the class $`Q`$ is parametrizing for $`S_a`$, and it follows that $`\varphi _a^{}(Q)`$ is parametrizing for $`S`$. We compute
$$\varphi _a^{}(Q)D=Q\varphi _{a}^{}{}_{}{}^{}(D)=Q(D_aaK_a)$$
$$=(C+P)(2aC+2aP+kP)=4a+k4a+2k.$$
If $`SF_n`$ and $`n>0`$, then we take $`Q:=C+nP`$. Then $`dim(C+nP)=n+1`$, and the image $`S_Q^{n+1}`$ of the associated map is a cone over a rational normal curve of degree $`n`$. By repeating projections from nonsingular points defined over $`𝕂`$, we obtain a birational map to the plane. Hence $`Q`$ is parametrizing, and therefore $`\varphi _a^{}(Q)`$ is parametrizing. We have
$$\varphi _a^{}(Q)D=(C+nP)(2aC+(an+k+2a)P)=an+2a+k.$$
Moreover, $`D`$ is nef, hence
$$0\varphi _a^{}(D)(C)=2a+kan,$$
hence $`\varphi _a^{}(Q)D4a+2k`$.
In case (b), we take $`Q:=P+E`$. The image of the associated map is a quadric in $`^3`$, hence $`Q`$ is parametrizing, as above. In this case, we get
$$\varphi _a^{}(Q)D=(P+E)(kP+2aP+aE)=4a+2k.$$
In case (c), we take $`Q=L`$, which is of course parametrizing. In this case, we get
$$\varphi _a^{}(Q)D=L(kPaK)=L(3aL+2kL(a+k)E)=3a+2k4a+2k.$$
Comprizing Theorems 6, 17, 19, and Remark 3, we can finally state:
###### Theorem 20.
For any rational surface $`S`$, the following bounds hold:
$$3\mathrm{level}(S)+\mathrm{keel}(S)pdeg(S)6\mathrm{level}(S)+2\mathrm{keel}(S).$$ |
warning/0506/quant-ph0506160.html | ar5iv | text | # MUTUAL INFORMATION OF BIPARTITE STATES AND QUANTUM DISCORD IN TERMS OF COHERENCE INFORMATION
## 1 Introduction
The investigation in this article is directed at the quantum correlations contained in a general, i. e., pure or mixed, bipartite state. By ”contained” is meant the von Neumann mutual information of the state. There are numerous other important approaches in the literature that are not limited to the mutual information . These will not be touched upon in this work.
No need to expand on the importance of this problem for quantum information theory, quantum communications, and quantum computers.
We will distinguish the two subsystems by $`1`$ and $`2`$. The former will be called ”the distant” subsystem, and the latter ”the nearby” one. We will distinguish ”local” properties of the nearby subsystem (or of the distant one), and ”global” ones of the bipartite state.
The approach of this article is based on the concept of coherence information. Coherence of an observable $`A`$ with respect to a quantum state $`\rho `$ and the incompatibility of the two have been simultaneously quantified by the concept of coherence information $`I_C(A,\rho )`$ . It is defined in three equivalent ways:
$$I_C(A,\rho )=S\left(\underset{l}{}P_l\rho P_l\right)S\left(\rho \right),$$
$`\left(1a\right)`$
where $`A=_la_lP_l`$ is the spectral form of the Hermitian operator $`A`$ in terms of distinct eigenvalues $`a_l`$, and $`S\left(\mathrm{}\right)`$ is the von Neumann entropy of a state. Further,
$$I_C(A,\rho )=S\left(\rho \right|\left|\underset{l}{}P_l\rho P_l\right),$$
$`\left(1b\right)`$
where $`S\left(\rho \right|\left|\sigma \right)`$ is the relative entropy, a known function of two states, and finally,
$$S\left(\rho \right)=S(A,\rho )+\underset{l}{}p_lS\left(P_l\rho P_l/p_l\right)I_C(A,\rho ),$$
$`\left(1c\right)`$
where $`S(A,\rho )=H\left(p_l\right)_lp_llogp_l`$ quantifies the uncertainty of $`A`$ in $`\rho `$ in terms of the Shannon entropy $`H\left(p_l\right)`$ of the probability distribution $`l:p_l\mathrm{tr}\left(P_l\rho \right).`$
The coherence information $`I_C(A,\rho )`$ quantifies also the quantumness in the relation between observable and state: The relation is quasi-classical if and only if $`A`$ and $`\rho `$ are compatible $`[A,\rho ]=0`$; in this and only in this case $`I_C(A,\rho )=0.`$
It will turn out that the coherence-information approach of this paper is closely connected with the Zurek concept of quantum discord. (It will be called shortly ”discord”.) Zurek introduced an approach in which the bipartite state $`\rho _{12}`$ is investigated by ”interrogating” it with a complete nearby subsystem observable $`A_2^c`$ , , , . The associated discord $`\delta _{A_2^c}\left(\rho _{12}\right)`$ appeared as the natural quantification of quantumness of the correlations. It is not entanglement. Also separable states, which, by definition of entanglement, do not contain it, are stated to have positive discord, showing quantumness in ”interrogation” by a concrete subsystem observable. Nevertheless, discord addresses, just like entanglement, though in a different way, the same basic problem of quantum correlations: What is there typically quantum mechanical in them?
In Zurek takes a thermodynamical approach to the study of the physical meaning of least discord $`\stackrel{ˇ}{\delta }inf_{\left\{A_2^c\right\}}\delta _{A_2^c}\left(\rho _{12}\right).`$ He does this using the idea of a quantum demon extracting locally work from $`\rho _{12}`$. He finds that $`\stackrel{ˇ}{\delta }`$ equals the (nonnegative) excess of work that a quantum demon can extract in comparison with a classical one. He also discusses how his approach relates to a similar thermodynamical approach of Oppenheim and the Horodecki family .
In a recent review article the Horodecki family, Oppenheim et al. gave a detailed presentation on ”local versus non-local information” . They discuss the connection between their approach and results with those of Zurek and his discord. Indirectly, the results of this article are connected also with this work.
Uhlmann gives an elementary presentation of an analogous approach to quantum correlations studies independently from both Zurek and the Horodecki school of thought .
It will be shown that discord is actually coherence-information excess (global minus local). This will make it possible to throw new light on the zero-discord problem. The ”interrogating” complete observable $`A_2^c`$ will be generalized to include also incomplete observables $`A_2`$. Then a string of relevant observables, each a function of the next, will be derived that will eliminate, what will be called, redundant noise, eliminate the garbled part of the information gain (on the state of the distant subsystem), and, finally, eliminate all quantumness - all this at the cost of diminishing the information gain.
The state $`\rho _L_lP_l\rho P_l`$ appearing in definitions (1a) and (1b) is the so-called Lüders mixture of $`\rho `$ with respect to $`A`$ , (relation (14.16) on p. 225 there). It is the non-selective (or entire-ensemble) version. (Some authors call it ”dephasing operation, e. g., .) The admixed Lüders state is $`P_l\rho P_l/p_l,`$ where $`p_l\mathrm{tr}\left(\rho P_l\right)`$ is the corresponding probability. It appears in the selective (or definite-result) version (utilized in (1c) e. g.).
To avoid unnecessary repetitions in the exposition, the following will be understood throughout the article: the physical term will be given priority, like ”observable” instead of ”Hermitian operator”, ”state” instead of ”statistical operator”, ”mixture” instead of ”decomposition of a statistical operator” (into a finite or infinite convex combination of statistical operators), ”compatibility” instead of ”commutation”, etc.
Observables will be restricted to discrete ones, and as a rule, given in spectral form like $`A=_la_lP_l`$ with all eigenvalues $`a_l`$ distinct. This will always be tacitly accompanied by the completeness relation (decomposition of the identity) $`_lP_l=1.`$ The sum ”$`_l`$ is finite or infinite as the case may be. If the sum is necessarily restricted to be finite for some claim to be valid, then it will be written, e. g., like $`_{l=1}^m,`$ and it will be understood that $`m`$ is an integer. If the spectrum may be finite or infinite, we will write $`\{a_l:l\},`$ etc.; if it is necessarily finite, we will write $`\{a_l:l=1,2,\mathrm{},m\}`$. Complete observables $`A^c=_la_lll`$ are written with the suffix ”c”.
If the given state $`\rho `$ has an infinite-dimensional null space, then also observables $`A`$ that have a continuous part in their spectrum can be considered for coherence-information studies under the restriction that the subspace spanned by the eigen-subspaces of $`A`$ contains the range of $`\rho `$ .
Functions of an observable amount to coarsenings in the spectrum of the latter. We will prefer the term ”coarsening” because it has the simple opposite ”refinement”. Both are order relations like ”smaller or equal” and ”larger or equal respectively”.
When an observable $`A=_la_lP_l`$ and a state $`\rho `$ are given, we will speak of ”detectable” eigenvalues $`a_l`$ or index values $`l`$ meaning those for which $`p_l\mathrm{tr}\left(\rho P_l\right)>0`$. The spectrum $`\{a_l:l\},`$ and the set of the index values $`\{l:l\}`$ will always be understood to be connected by a fixed one-to-one map, enabling us to talk of ”corresponding” eigenvalue etc.
Mixtures like $`\rho =_kw_k\rho _k`$, finite or infinite, will be understood in a formal, not operational sense, i. e., they express the fact that one can write $`\rho `$ in that way. The statistical weights will be called only ”weights”; they can be positive or zero; in the latter case $`\rho _k`$ need not be defined, nevertheless by definition $`w_k\rho _k=0`$ (and analogously for other entities than $`\rho _k`$). The states $`\rho _k`$ will be referred to as ”admixed states”.
A mixture $`\rho =_kw_k\rho _k`$ is orthogonal if $`kk^{}\rho _k\rho _k^{}=0.`$ An example is the Lüders mixture $`\rho _L=_lp_l\left(P_l\rho P_l/p_l\right).`$ Then the mixing property of entropy is valid: $`S\left(\rho _L\right)=H\left(p_l\right)+_lp_lS\left(P_l\rho P_l/p_l\right)`$ (see p. 242 in ).
Both for mixtures and for observables the subsystem will be exhibited in the index, e. g., $`_kw_k\rho _2^k,_la_lP_2^l`$.
Both in mixtures and in observables we will deal with coarsenings, and binary relations ”linked” and ”chained”. To distinguish the two cases, we will occasionally use the terms ”m-coarsening”, ”m-linked” and ”m-chained” for mixtures, and ”o-coarsening”, ”o-linked”, and ”o-chained” for observables.
Mentioning subsystems, we will often omit ”nearby”, and only say ”subsystem”. One should note that every general statement is symmetrical in the sense that one can interchange $`1`$ and $`2`$: the claim is either unchanged or one obtains the symmetrical equally valid claim. Thus, the stated choice of nearby and distant is arbitrary.
## 2 Role of Coherence Information in Bipartite Quantum Correlations
We take a bipartite state $`\rho _{12}`$ with its reductions $`\rho _s\mathrm{tr}_s^{}\rho _{12},s,s^{}=1,2ss^{},`$ and a subsystem observable $`A_2=_la_lP_2^l.`$ Two coherence informations $`I_C(A_2,\rho _2)`$ and $`I_C(A_2,\rho _{12})`$ appear. Also two Lüders mixtures $`\rho _2^L_lP_2^l\rho _2P_2^l`$ and $`\rho _{12}^L_lP_2^l\rho _{12}P_2^l`$ enter the scene. (Here $`P_2^l`$ is short for $`(1P_2^l).)`$ We utilize the notation:
$$l:p_l\mathrm{tr}\left(\rho _{12}P_2^l\right),$$
$`\left(2a\right)`$
$$l,p_l>0:\rho _{12}^lP_2^l\rho _{12}P_2^l/p_l,$$
$`\left(2b\right)`$
$$l,p_l>0:\rho _s^l\mathrm{tr}_s^{}\left(\rho _{12}^l\right),s,s^{}=1,2ss^{}.$$
$`\left(2c\right)`$
Next, we’ll need the entropy additivity accompanying any mixture $`\rho =_kw_k\rho _k`$:
$$S\left(\rho \right)=J+\underset{k}{}\left(w_kS\left(\rho _k\right)\right),$$
$`\left(3a\right)`$
$$J=\underset{k}{}\left(w_kS\left(\rho _k\right|\left|\rho \right)\right).$$
$`\left(3b\right)`$
(If proof is wanted for the known relation (3b), see proof of Lemma 4 in .) If the mixture is orthogonal, then $`J`$ takes the special form of the Shannon entropy $`H\left(w_k\right)`$ due to the mixing property. (See proposition 7 below for more on $`J`$.)
Now we consider a relevant decomposition of the mutual information $`I_{12}I\left(\rho _{12}\right)S_1+S_2S_{12}`$, where $`S_sS\left(\rho _s\right),s=1,2,12`$.
Theorem 1: A) The mutual information $`I_{12}`$ of any bipartite state $`\rho _{12}`$, when viewed in relation to any given discrete second-subsystem observable $`A_2`$, can be decomposed as follows:
$$I_{12}=J_{A_2}+\left(I_C(A_2,\rho _{12})I_C(A_2,\rho _2)\right)+\underset{l}{}p_lI\left(\rho _{12}^l\right),$$
$`\left(4a\right)`$
where
$$J_{A_2}\underset{l}{}p_lS\left(\rho _1^l\right|\left|\rho _1\right),$$
$`\left(4b\right)`$
and
$$\rho _1=\underset{l}{}p_l\rho _1^l$$
$`\left(4c\right)`$
is the distant mixture induced by $`A_2`$.
B) Each of the three terms on the RHS of (4a) is always nonnegative.
Proof: A) We utilize the entropy decompositions (3a) for $`\rho _1`$ and (1c) for $`\rho _s,s=2,12`$:
$$I_{12}S_1+S_2S_{12}=\left[J_{A_2}+\underset{l}{}p_lS\left(\rho _1^l\right)\right]+\left[H\left(p_l\right)+\underset{l}{}p_lS\left(\rho _2^l\right)I_C(A_2,\rho _2)\right]$$
$$\left[H\left(p_l\right)+\underset{l}{}p_lS\left(\rho _{12}^l\right)I_C(A_2,\rho _{12})\right]=RHS\left(4a\right).$$
This completes the proof of part A).
B) The first and the third terms on the RHS of (4) are obviously nonnegative. To prove that also the second term is nonnegative we need two auxiliary claims.
Corollary 1: Decomposition (4a) in application to the Lüders mixture $`\rho _{12}^L_lP_2^l\rho _{12}P_2^l`$ gives:
$$I\left(\rho _{12}^L\right)=J_{A_2}+\underset{l}{}p_lI\left(\rho _{12}^l\right).$$
$`\left(5\right)`$
Proof: Straightforward evaluation gives $`I_C(A_2,\rho _{12}^L)=I_C(A_2,\rho _2^L)`$ (or see proposition 1 below). $`\mathrm{}`$
Lemma 1: The inequality $`I\left(\rho _{12}^L\right)I\left(\rho _{12}\right)`$ is always valid.
Proof: As it is well known, the mutual information can be written in the form of relative entropy $`I_{12}=S\left(\rho _{12}\right||\rho _1\rho _2)`$. By this same formula also $`I\left(\rho _{12}^L\right)=S\left(_lP_2^l\rho _{12}P_2^l\right||\rho _1\left(_lP_2^l\rho _2P_2^l\right))`$. It was proved by Lindblad for the finite-dimensional case (Theorem on p. 149 there) that $`S\left(\mathrm{\Phi }\rho \right|\left|\mathrm{\Phi }\sigma \right)S\left(\rho \right|\left|\sigma \right)`$ for every two states $`\rho `$ and $`\sigma `$ and every completely positive trace-preserving map $`\mathrm{\Phi }`$. The inequality was extended to the infinite-dimensional case by Uhlmann . (It is unjustly called a theorem of Uhlmann instead of one of Lindblad and Uhlmann.)
Since $`\mathrm{\Phi }_lP_2^l\mathrm{}P_2^l`$ is such a map, the lemma is proved. $`\mathrm{}`$
End of proof of part B) of Theorem 1: Comparing (4a) and (5) and taking into account Lemma 1, one obtains
$$I\left(\rho _{12}\right)I\left(\rho _{12}^L\right)=I_C(A_2,\rho _{12})I_C(A_2,\rho _2)0$$
in the general case. $`\mathrm{}`$
In the classical discrete case one has a relation analogous to (3a) and (3b), and one analogous to (1c), but the latter with $`I_C`$ missing. Then a relation analogous to (4a) is obtained (analogously as in the proof of theorem 1), but without the excess of coherence information (the second term) on the RHS. Following Zurek , this term quantifies the quantumness in the mutual information and is called the quantum discord with respect to a complete or incomplete second-subsystem discrete observable $`A_2`$, and it is denoted by $`\delta _{A_2}\left(\rho _{12}\right)`$.
The following physical interpretation of (4a) suggests itself. The observable $`A_2`$ is a probe (or an ”interrogation”, cf ) into the quantum correlations in $`\rho _{12}`$ making subsystem $`2`$ the nearby one (the instrument measuring $`A_2`$ interacts directly with it), and subsystem $`1`$ the distant one (no interaction with the measuring apparatus). Applying (3a) and (3b) to the mixture (4c), one obtains
$$S\left(\rho _1\right)=\underset{l}{}\left(p_lS\left(\rho _1^l\right|\left|\rho _1\right)\right)+\underset{l}{}\left(p_lS\left(\rho _1^l\right)\right).$$
$`\left(6\right)`$
In view of (6), the first term on the RHS of (4a) is obviously the information gain about the distant subsystem acquired by the probe (cf , , ). The detectable eigenvalues $`a_l`$ of $`A_2`$ distinguish and enumerate the admixed states $`\rho _1^l`$, and the acquired information is the gain in the distant mixture (4c).
The third term on the RHS of (4a) is the amount of quantum correlations in $`\rho _{12}`$ that is inaccessible by the probe used. (As easily seen, it is zero if $`A_2`$ is complete.) We shall call it residual correlations. Both the first and the third term are entropy terms, i. e., as easily seen, they are concave with respect to mixtures. But since the mutual information appears with a minus sign in the subsystem entropy decomposition $`S\left(\rho _{12}\right)=S\left(\rho _1\right)I\left(\rho _{12}\right)+S\left(\rho _2\right),`$ the mentioned terms are actually convex as information quantities should be.
Discord $`\delta _{A_2}\left(\rho _{12}\right),`$ being, in general, excess coherence information, i. e., a difference of two information quantities:
$$\delta _{A_2}\left(\rho _{12}\right)=\left[I_C(A_2,\rho _{12})I_C(A_2,\rho _2)\right],$$
$`\left(7\right)`$
is neither convex nor concave (because coherence information is convex, cf proposition 5 in ). This fact gives some insight into Lieb’s result that mutual information is neither convex nor concave in the general case . Some authors call $`I_{12}`$ ”mutual entropy”. Having its behavior under mixing in view, it is neither information nor entropy. (See \- subsection III.c there - for a different point of view.)
Discord is a necessary accompaniment of the described probing into $`\rho _{12}`$ by $`A_2`$. It is due to quantumness of the correlations.
Assuming that the observable $`A_2=_la_lP_2^l`$ is incomplete, one may wonder how the terms in (4a) behave when $`A_2`$ is refined (down to a complete observable or just to a more complete one). By refinement is meant another observable
$$A_2^{}=\underset{l,q}{}a_{l,q}P_2^{l,q}$$
$`\left(8a\right)`$
(the range of $`q`$ depends on the value of $`l`$; for simplicity, this is omitted in notation). It is by definition such that it further decomposes the eigenprojectors of $`A`$, i. e.,
$$l:P_2^l=\underset{q}{}P_2^{l,q}.$$
This is refinement in an absolute sense, i. e., it does not depend on any state $`\rho _2`$. We need the generalization of this notion to state-dependent refinement.
Let besides $`A_2^{}`$ (cf (8a)) also $`A_2`$ and $`\rho _2`$ be given. Let $`l^{}`$ enumerate the detectable and $`l^{\prime \prime }`$ the undetectable eigenvalues of $`A_2`$ in $`\rho _2`$. Then
$$A_2=\underset{l^{}}{}a_l^{}P_2^l^{}+\underset{l^{\prime \prime }}{}a_{l^{\prime \prime }}P_2^{l^{\prime \prime }}.$$
$`\left(8b\right)`$
If
$$l^{}:P_2^l^{}=\underset{q}{}P_2^{l^{},q},$$
$`\left(8c\right)`$
then we say that $`A_2^{}`$ is a (state-dependent) refinement of $`A_2`$ with respect to $`\rho _2`$, and we write $`A_2^{}\stackrel{\rho _2}{}A_2`$. (The symbol ”$`\stackrel{\rho _2}{}`$” is to remind us that we are dealing with a reflexive and transitive binary relation - like ”larger or equal” - that is state dependent.)
Theorem 2: In refinement of $`A_2`$ by $`A_2^{}`$ with respect to $`\rho _2`$ (cf (8a)-(8c)), the reduction of a given arbitrary bipartite state $`\rho _{12}`$, the information gain and the discord remain equal or become larger, and the residual correlations remain the same or become smaller. To be explicit quantitatively, one can write (4a) and (4b) with respect to $`A_2^{}`$ as a two-step expression (as if the probing took place first with $`A_2`$, and then it was continued to $`A_2^{}`$):
$$I_{12}=\{\underset{l}{}\left(p_lS\left(\rho _1^l\right|\left|\rho _1\right)\right)+\underset{l,q}{}\left(p_lp_{l,q}S\left(\rho _1^{l,q}\right|\left|\rho _1^l\right)\right)\}+$$
$$\left\{\delta _{A_2}\left(\rho _{12}\right)+\underset{l}{}p_l\delta _{A_2^{}}\left(\rho _{12}^l\right)\right\}+\left\{\underset{l,q}{}\left(p_lp_{l,q}I\left(\rho _{12}^{l,q}\right)\right)\right\},$$
$`\left(9\right)`$
where the expressions in the large brackets are the information gain, the discord and the residual correlations respectively (and $`p_{l,q}\mathrm{tr}\left(P_2^{l,q}\rho _{12}^l\right)`$).
Proof is given in Appendix A.
Information gain is the basic purpose of the probe, hence, one wants it to be as large as possible. This is the reason why most studies are restricted to complete observables $`A_2^c`$. Then, whenever $`p_l>0,`$ the state $`l_2l_2\rho _2l_2l_2/p_l=l_2l_2`$ is pure, $`\rho _{12}^l=\rho _1^ll_2l_2`$ is uncorrelated, and $`I\left(\rho _{12}^l\right)=0.`$ Then (4a) is simplified to become
$$I\left(\rho _{12}\right)=J_{A_2^c}+\left(I_C(A_2^c,\rho _{12})I_C(A_2^c,\rho _2)\right).$$
$`\left(10\right)`$
It was argued in that taking the infimum of the discords in (10) (cf (7))
$$\stackrel{ˇ}{\delta }\left(\rho _{12}\right)inf_{\left\{A_2^c\right\}}\delta _{A_2^c}\left(\rho _{12}\right)$$
$`\left(11\right)`$
one may obtain an observable-independent quantum measure of quantumness in $`I_{12}`$. Vedral et al. take into account also generalized observables, and then, taking the supremum of the $`J_{A_2^c}=_lp_lS\left(\rho _1^l\right|\left|\rho _1\right)`$ expressions, they define the classical part of $`I_{12}`$.
## 3 On Zero Discord
As it is obvious from (7), a discord $`\delta _{A_2}\left(\rho _{12}\right)`$ can be zero either if both coherence informations are zero, then we call it strong zero, or if both coherence informations are positive but equal. We call this case weak zero.
A detailed analysis including open problems (at least for the author) on unachieved results is now presented.
### 3.1 Strong zero discord with an incomplete or complete observable
Proposition 1: Each of the following two relations is a necessary and sufficient condition for an observable $`A_2=_la_lP_2^l`$ to have a strong zero discord in a given bipartite state $`\rho _{12}`$:
$$[A_2,\rho _{12}]=0,$$
$`\left(12\right)`$
$$\rho _{12}=\underset{l}{}P_2^l\rho _{12}P_2^l.$$
$`\left(13\right)`$
Proof: Upon partial trace over the first subsystem, the commutation (12) becomes $`[A_2,\rho _2]=0.`$ Hence the sufficiency and the necessity of this condition is obvious.
Relation (12) is equivalent to
$$l:[P_2^l,\rho _{12}]=0.$$
$`\left(14\right)`$
The identity $`\rho _{12}=\left(_lP_2^l\right)\rho _{12},`$ idempotency and commutation then give (13). Conversely, (13) implies (14).
$`\mathrm{}`$
Remark 1: Relation (12) implies the local necessary condition $`[A_2,\rho _2]=0`$ for strong zero discord. A local sufficient condition is not possible in a nontrivial way. Namely, if such a condition were given in terms of $`A_2`$ and $`\rho _2`$ only, one could make the so-called purification: $`\rho _{12}\mathrm{\Psi }_{12}\mathrm{\Psi }_{12}`$ with $`\mathrm{tr}_1\rho _{12}=\rho _2`$ (the given local state). Then, relation (14) would imply, as easily seen,
$$\overline{l}:l:\left(1P_2^l\right)\mathrm{\Psi }_{12}=\delta _{l,\overline{l}}\mathrm{\Psi }_{12},$$
and further
$$l:P_2^l\rho _2=\delta _{l,\overline{l}}\rho _2.$$
This gives zero discord, but it also gives zero information gain $`J=0`$ because it does not decompose $`\rho _1`$ at all, and thus it is a trivial probe.
One wants to know what kind of state $`\rho _{12}`$ has a strong zero discord.
Definition 1: If a bipartite state $`\rho _{12}`$ is a nontrivial mixture of admixed states $`\rho _{12}^k`$
$$\rho _{12}=\underset{k}{}w_k\rho _{12}^k$$
$`\left(15a\right)`$
(all weights $`w_k`$ being positive) so that
$$kk^{}\rho _2^k\rho _2^k^{}=0,$$
$`\left(15b\right)`$
where $`k:\rho _2^k\mathrm{tr}_1\rho _{12}^k,`$ then $`\rho _{12}`$ is said to be mono-orthogonal (in the second subsystem).
Proposition 2: A bipartite state $`\rho _{12}`$ has a strong zero discord if and only if it is mono-orthogonal (in the second subsystem).
Proof: Sufficiency. Let us assume that a state $`\rho _{12}`$ for which (15a) and (15b) are valid is given. Let us, further, for each $`k`$ value denote by $`Q_2^k`$ the range-projector of $`\rho _2^k`$. Finally, let us define $`A_2_ka_kQ_2^k`$ with arbitrary but distinct eigenvalues $`a_k`$. Then one has $`k:\rho _{12}^k=Q_2^k\rho _{12}^kQ_2^k`$ (This is a known but not well known general relation. For proof cf relation (12a) in .) Hence (14) (changing what has to be changed) holds true.
Necessity. If $`\rho _{12}`$ has a strong zero discord with respect to an observable $`A_2=_la_lP_2^l,`$ then, according to the necessary condition (13), one can write $`\rho _{12}=_l^{}p_l\rho _{12}^l,`$ where the prim on the sum denotes that all $`\left(p_l=0\right)`$-terms are omitted, and $`l,p_l>0:\rho _{12}^lP_2^l\rho _{12}P_2^l/p_l`$. This is of the form (15a). Further, $`l,p_l>0:\rho _2^l\mathrm{tr}_1\rho _{12}^l=P_2^l\rho _2P_2^l/p_l,`$ and requirement (15b) (with $`l`$ instead of $`k`$) is obviously satisfied. $`\mathrm{}`$
Remark 2: Let it be locally known that $`\rho _{12}`$ is mono-orthogonal. This means that besides $`\rho _2`$ also an orthogonal projector decomposition $`_kQ_2^k=Q_2`$ of the range projector $`Q_2`$ of $`\rho _2`$ is given and it is known that it is associated with mono-orthogonality, i. e., $`Q_2^k`$ is the range projector of $`\rho _2^k\mathrm{tr}_1\rho _{12}^k,`$ where $`\rho _{12}^k`$ are the admixed mono-orthogonal states in (15a). Then, as easily seen, a local sufficient condition for strong zero discord is that each eigenprojector $`P_2^l`$ of $`A_2`$ be a sum of $`Q_2^k`$ projectors. This implies the necessary condition $`[A_2,\rho _2]=0`$ (because the $`Q_2^k`$ projectors commute with $`\rho _2`$). Nevertheless, it is not a necessary and sufficient condition, because it may require too much. A necessary and sufficient local condition cannot be given in view of lack of knowledge of the admixed mono-orthogonal states $`\rho _{12}^k`$ (cf remark 1).
### 3.2 Strong zero discord with a complete observable
The necessary and sufficient condition (12) is unchanged, but, since now $`A_2=_la_ll_2l_2,`$ (13) and (14) take the respective forms:
$$\rho _{12}=\underset{l}{}l_2l_2\rho _{12}l_2l_2,$$
$`\left(16\right)`$
and
$$l:[l_2l_2,\rho _{12}]=0.$$
$`\left(17\right)`$
Condition (16) was highlighted in (in a less elaborate context, without distinguishing strong and weak zero discord).
Proposition 3: A bipartite state $`\rho _{12}`$ has a strong zero discord with respect to a complete observable $`A_2=_la_ll_2l_2`$ if and only if it is a mixture of the form
$$\rho _{12}=\underset{l}{}p_l\rho _1^ll_2l_2.$$
$`\left(18\right)`$
Proof: Sufficiency. If (18) is valid, then so is (16).
Necessity. Since $`l:l_2l_2\rho _{12}l_2l_2=\left(l_2\rho _{12}l_2\right)l_2l_2=p_l\rho _1^ll_2l_2`$ (cf (2b) and (2c) with $`l_2l_2`$ instead of $`P_l`$). Thus, (16) implies (18). $`\mathrm{}`$
Proposition 4: A bipartite state $`\rho _{12}`$ has a strong zero discord with respect to some complete observable $`A_2`$ if and only if the state is mono-orthogonal (cf (15a) and (15b)), and
$$k:\rho _{12}^k=\rho _1^k\rho _2^k,$$
i. e., if it is simultaneously also separable.
Proof: Sufficiency. Let (15a) and (15b) be given, and let $`\rho _{12}`$ be simultaneously also separable as stated. Substituting each $`\rho _2^k`$ by its spectral form in terms of eigen-ray-projectors, one obtains $`\rho _{12}`$ as a mixture of the form (18) (changing what has to be changed).
Necessity. The form (18) is mono-orthogonal and simultaneously separable. $`\mathrm{}`$
Proposition 5: Let $`\rho _{12}`$ be a mixture of the form
$$\rho _{12}=\underset{k}{}w_k\rho _1^k\rho _2^k$$
$`\left(19\right)`$
with the validity of (15b) (cf definition 1 and proposition 4). Then a local sufficient condition for $`A_2=_la_lP_2^l`$ to give a strong zero discord for $`\rho _{12}`$ is:
$$k:[A_2,\rho _2^k]=0.$$
$`\left(20\right)`$
Proof: It is obvious in (19) that, on account of (20), $`A_2`$ commutes with $`\rho _{12}`$ (cf proposition 1). $`\mathrm{}`$
### 3.3 Weak zero discord
We begin by two general results, which play an auxiliary role in this subsection.
Lemma 2: Let $`\rho `$ be a state and $`A=_la_lP_l`$ an observable. Let, further, $`_nP_n=1`$ be an (orthogonal projector) decomposition of the identity such that
$$n:[P_n,\rho ]=[P_n,A]=0.$$
$`\left(21\right)`$
Then the following statistical decomposition of the coherence information ensues:
$$I_C(A,\rho )=\underset{n}{}w_nI_C(A,P_n\rho /w_n),$$
$`\left(22\right)`$
where $`n:w_n\mathrm{tr}\left(\rho P_n\right)`$.
Proof: On account of (21), one has the mixture $`\rho =_nw_n\left(P_n\rho /w_n\right),`$ and, $`[P_l,P_n]=0.`$ Hence,
$$I_C(A,\rho )S\left(\underset{l}{}P_l\rho P_l\right)S\left(\rho \right)=$$
$$S\left(\underset{n}{}w_n\underset{l}{}P_l\left(P_n\rho /w_n\right)P_l\right)S\left(\underset{n}{}w_n\left(P_n\rho /w_n\right)\right)=$$
$$H\left(w_n\right)+\underset{n}{}w_nS\left(\underset{l}{}P_l\left(P_n\rho /w_n\right)P_l\right)$$
$$\left[H\left(w_n\right)+\underset{n}{}w_nS\left(P_n\rho /w_n\right)\right]=\underset{n}{}w_nI_C(A,P_n\rho /w_n).$$
The symbol $`H\left(w_n\right)`$ denotes the Shannon entropy $`\mathrm{tr}\left(w_nlogw_n\right),`$ and the mixing property of entropy has been made use of. $`\mathrm{}`$
Lemma 3: Let $`\rho _{12}`$ be a bipartite state and $`A_2=_la_lP_2^l`$ a subsystem observable. Besides, let $`_nP_2^n=1`$ be a subsystem (orthogonal projector) decomposition of the identity such that
$$n:[P_2^n,\rho _{12}]=0\text{and}[P_2^n,A_2]=0.$$
$`\left(23\right)`$
Then the following statistical decomposition of the discord is valid:
$$\delta _{A_2}\left(\rho _{12}\right)=\underset{n}{}w_n\delta _{A_2}\left(P_2^n\rho _{12}/w_n\right),$$
$`\left(24\right)`$
where the mixture $`\rho _{12}=_n\left[w_n\left(P_2^n\rho _{12}/w_n\right)\right]`$ is due to (23).
Proof: Taking the first-subsystem partial trace in the first commutation relation in (23), one obtains $`[P_2^n,\rho _2]=0.`$ Hence, according to (7) and lemma 2,
$$\delta _{A_2}\left(\rho _{12}\right)=I_C(A_2,\rho _{12})I_C(A_2,\rho _2)=\underset{n}{}w_n\delta _{A_2}\left(P_2^n\rho _{12}/w_n\right).$$
$`\mathrm{}`$
Proposition 6: A sufficient condition for a weak zero discord of $`A_2`$ in $`\rho _{12}`$ is the mixture (19) of the latter with (15b) valid, further,
$$k:[A_2,Q_2^k],$$
$`\left(25a\right)`$
where $`Q_2^k`$ is the range projector of $`\rho _2^k`$, and finally, for at least one detectable value $`\overline{k}`$ of $`k`$ one has
$$[A_2,\rho _2^{\overline{k}}]0.$$
$`\left(25b\right)`$
Proof: Since $`k:\rho _2^k=Q_2^k\rho _2^kQ_2^k,`$ and $`\rho _1^k\rho _2^k=Q_2^k\left(\rho _1^k\rho _2^k\right)Q_2^k,`$ the assumptions of lemma 3 are satisfied with the decomposition $`_kQ_2^k=1`$. (The null-space projector of $`\rho _2`$, if it is nonzero, is joined to the $`Q_2^k`$.) Hence, one can write
$$\delta _{A_2}\left(\rho _{12}\right)=\underset{k}{}w_k\delta _{A_2}\left(\rho _1^k\rho _2^k\right)=0,$$
because uncorrelated states have zero mutual information, and this is an upper bound for the (nonnegative) discord (cf (7) and (4a)).
On the other hand, also the assumptions of lemma 2 are satisfied. Thus
$$I_C(A_2,\rho _2)=\underset{k}{}w_kI_C(A_2,\rho _2^k)w_{\overline{k}}I_C(A_2,\rho _2^{\overline{k}})>0.$$
In view of (7), the zero discord must be weak. $`\mathrm{}`$
Remark 3: One would like to know if the condition in Proposition 6 is also necessary, or if some other at least partially local necessary and sufficient condition is valid.
As it is well known, in quantum mechanics, unlike in the classical discrete case, the von Neumann mutual information $`I_{12}`$ can exceed the subsystem entropies, actually $`I_{12}2min(S\left(\rho _1\right),S\left(\rho _2\right)).`$ Any correlated pure bipartite state is a good example, because, as it is also well known, there $`I_{12}=2S\left(\rho _1\right)=2S\left(\rho _2\right).`$
Substituting (7) in (4a) and utilizing (3a), (4a) implies for any complete subsystem observable $`A_2^c`$
$$\delta _{A_2^c}=\underset{l}{}p_lS\left(\rho _1^l\right)+\left(I_{12}S_1\right).$$
$`\left(26\right)`$
If $`I_{12}`$ exceeds $`S_1`$, then (26) gives rise to a lower bound
$$\delta _{A_2^c}\left(I_{12}S_1\right)>0.$$
$`\left(27\right)`$
Thus, for such typically quantum states $`\rho _{12}`$ no choice of $`A_2^c`$ can give zero discord.
Cerf and Adami introduced quantum conditional entropies $`S\left(1|2\right)`$ . One has $`S\left(1|2\right)=S_1I_{12}.`$ If (27) is valid, then $`S\left(1|2\right)<0.`$ It is what Adami and Cerf call ”supercorrelations” .
The opposite-sign entity $`S\left(1|2\right)E\left(12\right)`$ is called ”directed entanglement” by Devetak and Staples . Its properties are discussed and it is applied to quantum communication. The same entity was called ”coherent quantum information” (not to be confused with ”coherence information” of the present study) by Schumacher and Nielsen with analogous discussion and application.
Remark 4: One would like to know if there can be zero discord between the case of mono-orthogonal and the case of states for which (27) is valid. In other words, one wonders if for some separable but not mono-orthogonal states and for some nonseparable but states for which $`I_{12}S_1,`$ one can find a complete subsystem observable $`A_2^c`$ giving zero discord.
Remark 5: It is desirable to learn if in the definition of the least discord $`\stackrel{ˇ}{\delta }inf_{\left\{A_2^c\right\}}\delta _{A_2^c}\left(\rho _{12}\right)`$ one can replace ”inf” by ”min” or not. In other words, it might be that there exist states $`\rho _{12}`$ for which $`\stackrel{ˇ}{\delta }`$ is ”irrational” in the sense that it can be reached by no $`A_2^c`$, but it can be arbitrarily well approximated by some $`\delta _{A_2^c}`$. One wants to see such states if they exist, or to see a proof that they do not exist. This is, of course, important also for $`\stackrel{ˇ}{\delta }=0`$.
The investigation in this section reveals that there is a number of open problems about the zero discord (contrary to a false impression one might mistakenly get, e. g., from ).
## 4 String of Relevant Coarsenings
### 4.1 Elaborate subsystem entropy decomposition
When a bipartite state $`\rho _{12}`$ is given and a subsystem observable $`A_2`$ is selected, then the subsystem entropy decomposition
$$S_{12}=S_1I_{12}+S_2$$
$`\left(28a\right)`$
can be viewed in the more elaborate way
$$S_{12}=\left\{\underset{l}{}p_lS\left(\rho _1^l\right)+J_{A_2}\left(\rho _1\right)\right\}\left\{J_{A_2}\left(\rho _1\right)+\delta _{A_2}\left(\rho _{12}\right)+\underset{l}{}p_lI\left(\rho _{12}^l\right)\right\}+$$
$$\left\{H\left(p_l\right)I_C(A_2,\rho _2)+\underset{l}{}p_lS\left(\rho _2^l\right)\right\}$$
$`\left(28b\right)`$
(cf (2a)-(2c), (3a) and (3b), (4a), (7), and (1c)). Naturally, $`J_{A_2}\left(\rho _1\right)=J_{A_2}\left(\rho _{12}\right).`$ It is understood that each expression in large brackets in (28b) equals the corresponding term on the RHS of (28a).
The elaborate subsystem entropy decomposition (28b) can be interpreted physically as follows. The subsystem observable $`A_2`$ is chosen to ”interrogate” the uncertainty in the distant subsystem $`1`$; the measure of the latter is $`S_1`$. On account of this, $`S_1`$ is broken up into $`_lp_lS\left(\rho _1^l\right),`$ the part of $`S_1`$ that is inaccessible to our ”interrogation” (or the residual part), and $`J_{A_2}\left(\rho _1\right),`$ the information gain. The mutual information $`I_{12}`$, which quantifies the total quantum correlations in $`\rho _{12}`$, is decomposed into the mentioned information gain $`J_{A_2}\left(\rho _1\right),`$ the discord $`\delta _{A_2}\left(\rho _{12}\right),`$ and $`_lp_lI\left(\rho _{12}^l\right),`$ which is the part that is not made use of in the chosen ”interrogation” (residual correlations). The appearance of the information gain in $`I_{12}`$ shows that the quantum correlations in $`\rho _{12}`$ act as an information channel, transferring the information gain from subsystem $`1`$ to subsystem $`2`$. The discord appears because, unless $`A_2`$ is compatible with $`\rho _{12}`$, there is a part of the correlations that is unsuitable for the mentioned transfer of the information gain, which is a quasi-classical notion. This is why it is said that the discord quantifies the quantumness of the correlations (regarding $`A_2`$). Finally, the uncertainty in $`\rho _2`$, i. e., $`S_2`$ is broken up into $`H\left(p_l\right)_lp_llogp_l=S(A_2,\rho _2),`$ the entropy (or amount of uncertainty) of $`A_2`$ in the state of the second subsystem; into the coherence or incompatibility information $`I_C(A_2,\rho _2),`$ which is again a necessary accompaniment of our ”interrogation” due to the quantumness of $`\rho _2`$; and into $`_lp_lS\left(\rho _2^l\right),`$ which is the amount of uncertainty in $`\rho _2`$ inaccessible to $`A_2`$ (residual uncertainty).
It should be noted that (28b) does not describe a process; it only gives a relevant quantitative view of $`\rho _{12}`$. In other words, what the quantum correlations in $`\rho _{12}`$ do, among other things, is to transfer the information gain $`J_{A_2}\left(\rho _1\right)`$ from subsystem $`1`$ to subsystem $`2`$. Now it is natural to ask how we can extract it from subsystem $`2`$. Evidently, the thing to do is to measure $`A_2`$ on the nearby subsystem $`2`$, i. e., locally (see section VI). But then one extracts the amount of information $`H\left(p_l\right)`$, and not $`J_{A_2}\left(\rho _1\right)`$. This motivates the rest of investigation in this section.
### 4.2 Information gain $`J`$
It is the aim of this subsection to understand how the uncertainty $`H\left(p_l\right)=S(A_2,\rho _2)`$ and the information gain $`J_{A_2}\left(\rho _1\right)`$ relate to each other. We begin by a precise understanding of $`J_{A_2}\left(\rho _1\right).`$
Proposition 7: If $`\rho =_{l=1}^mp_l\rho ^l`$ is an arbitrary mixture of a finite number of admixed states, then (3a) and (3b) are valid. Besides,
$$0J\left(\rho \right)H\left(p_l\right),$$
$`\left(29\right)`$
and $`J\left(\rho \right)=0`$ if and only if $`l,p_l>0:\rho ^l=\rho `$ (total overlap), and $`J\left(\rho \right)=H\left(p_l\right)`$ if and only if $`\left(ll^{}\right),p_l>0<p_l^{}:\rho ^l\rho ^l^{}=0`$ (pairwise orthogonality).
Proof: The first inequality in (29) is obvious from (3b). The second one is proved in the review article of Wehrl (relation (2.3) there).
One has $`J=0`$ if and only if in (3b) (changing what has to be changed) $`p_l>0S\left(\rho ^l\right|\left|\rho \right)=0.`$ It is well known that relative entropy is zero if and only if the two states in it coincide.
It is standard knowledge that the so-called mixing property holds true: if the admixed states $`\rho ^l`$ are pairwise orthogonal, then $`J=H\left(p_l\right).`$ The converse statement, that $`J=H\left(p_l\right)`$ implies orthogonality of the $`\rho ^l`$, is not proved anywhere known to the author of this study. Therefore, its somewhat lengthy proof, through auxiliary lemmata, is given in Appendix B. $`\mathrm{}`$
The quantity $`H\left(p_l\right)`$ is called the mixing entropy of the mixture at issue. But it is only the upper possible extreme value of the information gain $`J`$. It is obvious from proposition 7 that the excess $`\left(H\left(p_l\right)J\right)`$ (or how much is missing in the information gain) quantifies the overlap of the admixed states. It is zero if and only if there is no overlap (the admixed states are orthogonal). It is maximal, i. e., equal to $`H\left(p_l\right),`$ in case of total overlap, when one is dealing only with an apparent mixture.
Remark 6: It is desirable to have the extension of proposition 7 to the case of infinitely many admixed states.
To clarify what is apparent and what is genuine in a mixture, we consider two trivial lemmata.
Lemma 4: Let us take a mixture
$$\rho =\underset{s}{}p_s\rho ^s,S\left(\rho \right)=\underset{s}{}p_sS\left(\rho ^s\right)+J,$$
$`\left(30a\right)`$
and a refinement of it
$$s,p_s>0:\rho ^s=\underset{k_s}{}w_{k_s}\rho ^{k_s},\rho =\underset{s}{}\underset{k_s}{}p_sw_{k_s}\rho ^{k_s}.$$
$`\left(30b\right)`$
Then the residual entropy is non-increasing, whereas the information gain and the mixing entropy are non-decreasing. More precisely (in obvious notation):
$$S\left(\rho \right)=\underset{s}{}\underset{k_s}{}p_sw_{k_s}S\left(\rho ^{k_s}\right)+\left\{\underset{s}{}\left(p_sJ^s\right)+J\right\},$$
$`\left(31a\right)`$
$$H\left(p_sw_{k_s}\right)=H\left(p_s\right)+\underset{s}{}p_sH\left(w_{k_s}\right).$$
$`\left(31b\right)`$
Proof is straightforward.
Lemma 5: If the refinement in a mixture is done through mere repetition, i. e., if $`s,p_s>0:k_sk_s^{}\rho ^{k_s}=\rho ^{k_s^{}},`$ then the residual entropy and the information gain remain the same.
Proof is obvious from (31a) if one takes into account that $`s,p_s>0:J^s=0`$. $`\mathrm{}`$
It is now seen that the information gain is insensitive to apparent mixing (or repetition of the admixed states); it depends only on the genuine mixing, i. e., on the distinct admixed states. Contrariwise, the mixing entropy is insensitive to the distinction between genuine and apparent mixing, i. e., it increases whenever at least one of the refined probability distributions is nontrivial. Therefore, in spite of the fact that $`\left(H\left(p_l\right)J\right)`$ does quantify the overlap in the given mixture, which may contain repetition of admixed states, it can be diminished on the basis of (31b).
### 4.3 Essential noise and garbled information
Definition 2: If a given mixture $`\rho =_lp_l\rho _l`$ is rewritten without repetition of the admixed states with the use of a new index $`s`$, the expression $`\left(H\left(p_s\right)J\right)`$ quantifies the essential overlap in the mixture, i. e., the one due to the genuine mixing of the distinct admixed states. The original quantity of overlap is the sum of the quantity of essential overlap and of that of redundant overlap: $`\left(H\left(p_l\right)J\right)=\left(H\left(p_s\right)J\right)+\left(H\left(p_l\right)H\left(p_s\right)\right)`$.
One can see in (31b) that $`\left(H\left(p_l\right)H\left(p_s\right)\right)`$ is the increase in the mixing entropy due to repetition of admixed states.
Returning to the elaborate subsystem entropy decomposition (28b), we see that at best we can extract the information gain $`H\left(p_l\right)`$ from subsystem $`2`$ by measuring the subsystem observable $`A_2`$ (which is simultaneously the measurement of $`\left(1A_2\right)`$ in the bipartite state $`\rho _{12}`$). The difference $`\left(H\left(p_l\right)J_{A_2}\right),`$ corresponding to the overlap in the distant mixture $`\rho _1=_lp_l\rho _1^l,`$ appears now as noise. In accordance with definition 2, this noise is the sum of an essential term and a redundant term. One cannot eliminate the former (without changing drastically $`A_2`$, i. e., without taking another subsystem observable that is not a function of $`A_2`$) because the essential term is due to the overlap of the distinct admixed states in $`\rho _1`$, but one can dispose of the redundant noise by sheer coarsening.
Theorem 3: There exists one and only one coarsening $`B_2^{ess}`$ of $`A_2`$ in which the redundant noise is and the essential noise is not eliminated, and the induced distant mixture $`\rho _1=_sp_s\rho _1^s`$ is equal to the one obtained due to $`A_2`$ but rewritten with positive weights and without repetitions in the admixed states. To obtain the subsystem observable $`B_2^{ess}`$, one defines the following equivalence relation in the detectable spectrum of $`A_2`$: $`ll^{}`$ if $`\rho _1^l=\rho _1^l^{}`$ (cf (2c)). Further, enumerating by $`s`$ the obtained equivalence classes $`\{𝒞_s:s\},`$ one defines
$$B_2^{ess}\underset{s}{}b_sP_2^s,$$
$`\left(32a\right)`$
where $`\{b_s:s\}`$ is an arbitrary set of distinct nonzero real numbers, and
$$s:P_2^s\underset{l𝒞_s}{}P_2^l.$$
$`\left(32b\right)`$
Proof: Since $`s:p_s\mathrm{tr}\left(\rho _{12}P_2^s\right)=\left(_{l𝒞_s}p_l\right)>0,`$ and
$$\rho _1^sp_s^1\mathrm{tr}_2\left(\rho _{12}P_2^s\right)=\underset{l𝒞_s}{}\left(p_l/p_s\right)p_l^1\mathrm{tr}_2\left(\rho _{12}P_2^l\right)=$$
$$\underset{l𝒞_s}{}\left(p_l/p_s\right)\rho _1^l=\rho _1^{\overline{l}},$$
where $`\overline{l}`$ is any index value from the class $`𝒞_s`$. Thus, $`B_2^{ess}`$ does induce the desired mixture for $`\rho _1`$. It is evidently the unique coarsening of $`A_2`$ doing this because every coarsening has to break up the detectable spectrum of $`A_2`$ into classes, and the desired purpose cannot be achieved in any other way. $`\mathrm{}`$
In general, the information gain $`J`$ is garbled because in the measurement of $`A_2`$ it appears necessarily with (inseparable) essential noise $`\left(H\left(p_s\right)J\right)`$. (For a precise definition of ”garbled information gain” see the last but one term in (38) below.)
Needles to say that the expounded procedure of eliminating redundant noise is analogous in the classical discrete case of probability distributions.
### 4.4 Orthogonal distant mixture, pure information gain and twin observables
As it is obvious from proposition 7 and (28b), if the distant mixture $`\rho _1=_{l=1}^mp_l\rho _1^l`$ is orthogonal, and only in this case, the essential noise is zero. Then, one has pure information: $`J_{A_2}=H\left(p_l\right)=S(A_2,\rho _2).`$ In this case there is no redundant noise either. It may happen that orthogonality is achieved only after disposing of the redundant noise. Therefore, we concentrate on $`B_2^{ess}=_sb_sP_2^s`$ and the corresponding distant mixture $`\rho _2=_sp_s\rho _2^s`$, but to make the results more general, the suffix ”ess” is omitted.
Let $`Q_1^s`$ be the range projector of $`\rho _1^s`$. Orthogonality of the above mixture amounts to
$$Q_1^sQ_1^s^{}=\delta _{s,s^{}}Q_1^s,$$
$`\left(33a\right)`$
and one has
$$\underset{s}{}Q_1^s=Q_1,$$
$`\left(33b\right)`$
where $`Q_1`$ is the range projector of the distant state $`\rho _1`$. In this case, we prove the following result.
Proposition 8: Assuming positivity of all the probabilities $`p_s`$ and the validity of $`\left(_sP_2^s\right)\rho _2=\rho _2,`$ if the distant mixture $`\rho _1=_sp_s\rho _1^s`$ (cf (2a)-(2c) changing what has to be changed) is orthogonal, then
$$\left(\underset{s}{}P_2^s\right)\rho _{12}=\rho _{12}=\left(\underset{s}{}Q_1^s\right)\rho _{12},$$
$`\left(34a\right)`$
and
$$s:Q_1^s\rho _{12}=P_2^s\rho _{12}$$
$`\left(34b\right)`$
are valid.
Proof: Let $`Q_2`$ be the range projector of the nearby state $`\rho _2`$. The relation $`\left(_sP_2^s\right)\rho _2=\rho _2`$ then implies $`\left(_sP_2^s\right)Q_2=Q_2`$ (see Appendix A in ). Since one can always write $`\rho _{12}=Q_2\rho _{12}`$ (cf relation (12a) in ), the first equality in (34a) follows.
The relation (33b) and the fact that one can write $`\rho _{12}=Q_1\rho _{12},`$ then make also the second equality in (34a) seen to be valid.
Next, we prove that
$$ss^{},Q_1^sP_2^s^{}\rho _{12}=0.$$
$`\left(35\right)`$
For unequal index values one has $`\mathrm{tr}\left(Q_1^sP_2^s^{}\rho _{12}\right)=p_s\mathrm{tr}\left(Q_1^s\rho _1^s^{}\right)=p_s\mathrm{tr}\left(Q_1^s\left(Q_1^s^{}\rho _1^s^{}\right)\right)=0.`$ Further, $`0=\mathrm{tr}\left(Q_1^sP_2^s^{}\rho _{12}\right)=\mathrm{tr}\left(\left(Q_1^sP_2^s^{}\right)\rho _{12}\left(Q_1^sP_2^s^{}\right)\right),`$ and $`\left(Q_1^sP_2^s^{}\right)\rho _{12}\left(Q_1^sP_2^s^{}\right)=0`$ is well known to ensue. Then, the Lemma of Lüders ( or see FN 16 in ) entails the claimed relation (35).
Finally, utilizing relations (34a) and (35), one can argue as follows. $`Q_1^s\rho _{12}=Q_1^s\left(_s^{}P_2^s^{}\right)\rho _{12}=Q_1^sP_2^s\rho _{12}=P_2^s\left(_s^{}Q_1^s^{}\right)\rho _{12}=P_2^s\rho _{12}`$ as claimed in (34b). $`\mathrm{}`$
If one defines a first-subsystem observable $`B_1_sb_sQ_1^s`$ with arbitrary but distinct nonzero detectable eigenvalues $`\{b_s:s\},`$ then, according to Theorem 1 in and the theorem on so-called twin observables (p. 052321-3 in ) imply that proposition 8, actually, gives one more necessary and sufficient condition for $`(B_1,B_2)`$ to be twin observables in $`\rho _{12}`$.
Twin observables have a number of remarkable properties (cf also and the references therein). For this study an important implication is that $`[B_i,\rho _i]=0,i=1,2`$ (cf the mentioned Theorem 1 in ).
Two obvious consequences on the elaborate subsystem entropy decomposition (28b), which is the basic object of this study, follow:
$$I_C(B_2,\rho _2)=0=I_C(B_1,\rho _1),$$
$`\left(36a\right)`$
and, on account of (7), $`\delta _{A_2}\left(\rho _{12}\right)=I_C(B_2,\rho _{12}).`$
Thus, in this case (28b) simplifies to
$$S_{12}=S_1I_{12}+S_2=\left\{\underset{s}{}p_sS\left(\rho _1^s\right)+H\left(p_s\right)\right\}$$
$$\left\{H\left(p_s\right)+I_C(B_2,\rho _{12})+\underset{s}{}p_sI\left(\rho _{12}^s\right)\right\}+\left\{H\left(p_s\right)+\underset{s}{}p_sS\left(\rho _2^s\right)\right\},$$
where the mixing property is utilized for the orthogonal mixture $`\rho _1=_sp_s\rho _s`$.
If $`I_C(B_2,\rho _{12})>0,`$ then we have the case of so-called correlations incompatibility (cf Section 6 in ), in which the discord equals the coherence or incompatibility information of $`B_2`$ in $`\rho _{12}`$. Besides, there is no quantumness in $`\rho _2`$ with respect to $`B_2`$. (One has global coherence without local coherence.)
The quantity of uncertainty $`S\left(\rho _2\right)`$ of the nearby subsystem state now (possibly) exceeds the quantity of uncertainty $`S(B_2,\rho _2)`$ of the obsevable $`B_2`$ in $`\rho _2`$, which equals the pure information gain $`H\left(p_s\right)=J_{B_2}.`$
The assumption $`\left(_sP_2^s\right)\rho _2=\rho _2`$ is satisfied for $`B_2^{ess}`$ due to the very definition of the indices $`s`$ (all detectable $`l`$ values of $`A_2`$ are used up in it). Besides, on account of the definition of $`B_2^{ess}`$, all probabilities $`p_s`$ are positive.
So far in this subsection we had in mind the special case when the distant mixture $`\rho _1=_sp_s\rho _1^s`$ without repetition in the admixed states turns out orthogonal. Now we return to the general case and prove that there always exists a (possibly trivial) unique minimal coarsening $`C_2=_tc_tP_2^t`$ of $`B_2^{ess}`$, and, by consequence, of $`A_2`$, that gives an orthogonal distant mixture and, by a definition analogous to the above of $`B_1`$, an observable $`C_1=_tc_t^{}Q_1^t`$ that is its twin observable.
### 4.5 Minimal orthogonal coarsening of a mixture
Before we proceed, we first expound some relevant properties of mixtures as far as orthogonal coarsenings of them are concerned.
Lemma 6: For any two states $`\rho `$ and $`\rho ^{}`$ one has $`\mathrm{tr}\left(\rho \rho ^{}\right)0,`$ and $`\mathrm{tr}\left(\rho \rho ^{}\right)=0`$ if and only if $`\rho \rho ^{}=0`$.
Proof: Always $`\mathrm{tr}\left(\rho \rho ^{}\right)=\mathrm{tr}\left(\rho ^{1/2}\rho ^{}\rho ^{1/2}\right)0`$ because $`\rho ^{1/2}\rho ^{}\rho ^{1/2}`$ is a positive operator. Sufficiency of orthogonality for trace orthogonality is obvious. Necessity is seen as follows: $`\mathrm{tr}\left(\rho \rho ^{}\right)=0`$ implies $`\rho ^{1/2}\rho ^{}\rho ^{1/2}=0,`$ and this has, due to the Lemma of Lüders (), $`0=\rho ^{}\rho ^{1/2}=\rho ^{}\rho `$ as its consequence. $`\mathrm{}`$
Definition 3: Let $`\rho =_kw_k\rho _k`$ be a mixture with positive weights and without repetitions of the admixed states $`\rho _k`$. We say that the states $`\rho _k`$ and $`\rho _k^{}`$ are linked if $`\mathrm{tr}\left(\rho _k\rho _k^{}\right)>0.`$ If $`\rho _k`$ and $`\rho _k^{}`$ are such that there exists an integer $`n,n=1,\text{or}2,\text{or}\mathrm{},`$ and there can be found a chain of admixed states $`\{\rho _{k_i}:i=1,2,\mathrm{},n\}`$ such that $`\rho _k=\rho _{k_1},`$ $`\rho _{k_n}=\rho _k^{},`$ and any two neighboring states in the chain are linked, then we say that $`\rho _k`$ and $`\rho _k^{}`$ are chained, and we speak of m-chaining.
Definition 4: We say that a mixture $`\rho =_tw_t\rho _t`$ is a coarsening of another mixture $`\rho =_sw_s\rho _s,`$ the latter being without repetition in the admixed states and with positive weights, if the index set $`\{s:s\}`$ is partitioned into m-classes $`𝒯_t`$: $`\{s:s\}=_t𝒯_t`$ (the sum stands for the union of the non-overlapping classes), is enumerated by $`t`$, and $`\rho _t=_{s𝒯_t}\left(w_s/w_t\right)\rho _s,`$ where $`t:w_t\left(_{s𝒯_t}w_s\right)`$. In this case we speak of m-coarsening.
Proposition 9: Let $`\rho =_sp_s\rho _s`$ be a mixture with all weights positive and without repetition. Let, further, another mixture $`\rho =_tw_t\rho _t`$ be a coarsening of the former mixture, obtained by chaining (chained m-coarsening). Then the latter mixture is orthogonal, and it is minimal as such, i. e., if also $`\rho =_up_u\rho _u`$ is an orthogonal coarsening of the initial mixture, then it is also a coarsening of its chained m-coarsening.
Proof: Orthogonality can be proved as follows. Let $`tt^{},`$ and let $`s𝒯_t,`$ and $`s^{}𝒯_t^{}.`$ We assume ab contrario that $`\mathrm{tr}\left(\rho _s\rho _s^{}\right)>0.`$ Then, according to definitions 4 and 3, $`\rho _s`$ and $`\rho _s^{}`$ are linked, and hence belong to the same m-class $`𝒯_t`$ contrary to assumption. Hence, $`\rho ^s\rho ^s^{}=0`$ (cf lemma 6), implying $`\rho ^t\rho ^t^{}=0`$ (cf definition 4).
Minimality is proved in the following way. Let the partitioning $`\{s:s\}=_u𝒰_u`$ define an orthogonal coarsening $`\rho =_up_u\rho ^u`$ in analogy with definition 4. Considering the initial mixture $`\rho =_sp_s\rho ^s,`$ we assume that two distinct index values $`s,s^{}`$ are m-linked (cf definition 3). Lemma 6 claims that $`\rho ^s`$ and $`\rho ^s^{}`$ then cannot be orthogonal; hence $`s`$ and $`s^{}`$ must belong to one and the same m-class $`𝒰_u`$. Next, let $`s`$ and $`s^{}`$ be chained. Then any two neighboring index values in the chain belong to one and the same m-class $`𝒰_u`$, entailing the fact that also $`s`$ and $`s^{}`$ belong to the same m-class. Thus, any m-class $`𝒯_t`$ is a subset of some m-class $`𝒰_u`$. This means that the u-mixture is a coarsening of the t-mixture, and the latter is thus proved to be minimal. $`\mathrm{}`$
### 4.6 The pure part of information gain
We return now to our investigation of an arbitrary bipartite state $`\rho _{12}`$. We have defined $`B_2^{ess}=_sb_sP_2^s`$ to eliminate redundant noise.
Definition 5: We define $`C_2^{tw}_tc_tP_2^t`$ as a coarsening of $`B_2^{ess}`$ that induces m-chaining (cf definitions 4 and 3) of the distant mixture $`\rho _1=_sp_s\rho _1^s`$ (induced by $`B_2^{ess}`$), and by a spectrum $`\{c_t:t\}`$ consisting of any distinct nonzero real numbers.
If $`\{Q_1^t:t\}`$ are the range projectors of the distant admixed states $`\rho _1^t,`$ then defining, further, $`C_1^{tw}_tc_t^{}Q_1^t,`$ ( the eigenvalues $`c_t^{}`$ any distinct and nonzero real numbers), then, according to proposition 8 and the discussion after its proof, one obtains twin observables $`(C_1^{tw},C_2^{tw}).`$
Corollary 2: In case $`C_2^{tw}`$ is nontrivial, one has two parallel orthogonal mixtures with the common index $`t`$, the nearby one $`\rho _2=_tp_t\rho _2^t,`$ and the distant one $`\rho _1=_tp_t\rho _1^t.`$ In general, $`\rho _{12}`$ is not a mixture of the global states $`\rho _{12}^tP_2^t\rho _{12}P_2^t/p_t,`$ which give $`\rho _2^t`$ and $`\rho _1^t`$ as their reductions. The global states $`\rho _{12}^t`$ are biorthogonal, i. e., $`tt^{}\rho _i^t\rho _i^t^{}=0,i=1,2.`$
Since $`C_2^{tw}`$ is a coarsening of $`B_2^{ess}`$, the information gain $`J_{C_2^{tw}}`$ of the former is not larger than that of the latter (see theorem 2), i. e.,
$$J_{C_2^{tw}}J_{B_2^{ess}}=J_{A_2},$$
$`\left(37a\right)`$
and
$$J_{C_2^{tw}}=H\left(p_t\right)H\left(p_s\right)H\left(p_l\right)$$
$`\left(37b\right)`$
(cf (31b)). One should remember that $`H\left(p_s\right)J_{A_2},`$ due to (possible) essential noise.
Observable coarsening (or o-coarsening) ”$`\stackrel{\rho _2}{}`$” is the opposite relation to (state-dependent) observable refinement (or o-refinement) explained in section II. It is a reflexive and transitive binary relation, i. e., it is a partial order in the set of all observables. One has
$$C_2^{tw}\stackrel{\rho _2}{}B_2^{ess}\stackrel{\rho _2}{}A_2$$
$`\left(37c\right)`$
parallelling (37a) and (37b).
Returning to the elaborate subsystem entropy decomposition (28b), and having the relations (37a) and (37b) in mind, one can write
$$H\left(p_l\right)=S(A_2,\rho _2)=\left\{H\left(p_l\right)H\left(p_s\right)\right\}+$$
$$\left\{H\left(p_s\right)J_{B_2^{ess}}\right\}+\left\{J_{B_2^{ess}}H\left(p_t\right)\right\}+H\left(p_t\right).$$
$`\left(38\right)`$
The physical interpretation of (38) goes as follows. The entropy $`S(A_2,\rho _2)`$ (quantifying the uncertainty) of the initial subsystem observable $`A_2`$ in the nearby local state $`\rho _2`$ consist of redundant noise $`\left\{H\left(p_l\right)H\left(p_s\right)\right\},`$ of essential noise $`\left\{H\left(p_s\right)J_{B_2^{ess}}\right\},`$ of garbled information gain $`\left\{J_{B_2^{ess}}H\left(p_t\right)\right\},`$ and, finally of pure information gain $`H\left(p_t\right).`$ Naturally, each of the terms is positive or zero, as the case may be. The latter occurs when the corresponding subsystem observable is trivial, i. e., when it has only one detectable eigenvalue (the probability of which is then, of course, one).
For different choices of $`A_2`$ one may come to different $`C_2^{tw}`$. One may wonder if there always exists one $`C_2^{tw}`$ for all choices of $`A_2`$, as refined as possible. This is not true in the general case. Taking as an example the well-known singlet pure bipartite state, it is easy to see that for any choice of a nontrivial $`A_2`$, one has $`A_2=B_2^{ess}=C_2^{tw},`$ and one obtains thus infinitely many different $`C_2^{tw}`$ observables that are all complete.
The case of pure bipartite states deserves a separate discussion (see subsection H). But first we again need some more general theory.
### 4.7 Minimal compatible coarsening of an observable
We begin by some relevant theory on relation between observable and state.
Definition 6: We say that two (equal or distinct) index values $`t`$ and $`t^{}`$ of detectable eigenvalues $`c_t`$ and $`c_t^{}`$ respectively of a given observable $`C=_{t^{\prime \prime }}c_{t^{\prime \prime }}P_{t^{\prime \prime }}`$ are linked with respect to a given state $`\rho `$ if $`P_t\rho P_t^{}0.`$ When $`t`$ and $`t^{}`$ are such that there exists an integer $`n,n=1,\text{or}2,\text{or}\mathrm{},`$ and there can be found a chain of index values $`\{t_i:i=1,2,\mathrm{},n\}`$ such that $`t=t_1,`$ $`t_n=t^{},`$ and any two neighboring index values in the chain are linked, then we say that $`t`$ and $`t^{}`$ are chained. Occasionally, when it is desirable to make a distinction with respect to m-linking and m-chaining, we shall speak of o-linking (short for observable-linking) and o-chaining.
O-chaining includes o-linking, and it is reflexive, symmetric and transitive, i. e., an equivalence relation in the detectable part of the spectrum of the observable $`C`$.
Definition 7: Let $`C=_tc_tP_t`$ be a given observable, and let the observable $`D`$ be a coarsening of $`C`$ defined by means of chaining in the detectable part of the spectrum of the latter with respect to a given state $`\rho `$ (cf definition 6):
$$D\underset{k}{}d_kP_k,k:P_k\underset{t𝒞_k}{}P_t,$$
$`\left(39a\right)`$
where
$$\{t:t,p_t\mathrm{tr}\left(\rho P_t\right)>0\}=\underset{k}{}𝒞_k$$
$`\left(39b\right)`$
partitions the detectable part of the spectrum of $`C`$ into the equivalence classes $`𝒞_k`$ obtained by o-chaining, and the eigenvalues of $`D`$ are arbitrary distinct nonzero real numbers. We call $`D`$ the chained coarsening of $`C`$ with respect to $`\rho `$.
Proposition 10: The chained coarsening $`D\left(\stackrel{\rho }{}C\right)`$ given in definition 7 is compatible with $`\rho `$, $`[D,\rho ]=0`$. It is the most refined coarsening of $`C=_tc_tP_t`$ compatible with $`\rho `$, i. e., if $`\overline{C}_j\overline{c}_j\overline{P}_j,`$ is a coarsening of $`C`$, $`j:\overline{P}_j_{t\overline{𝒞}_j}P_t`$ with arbitrary distinct nonzero real eigenvalues of $`\overline{C}`$ and $`\{t:t,p_t>0\}=_j\overline{𝒞}_j`$ a partitioning of the detectable part of the spectrum of $`C`$, such that $`\overline{C}`$ is compatible with $`\rho `$, $`\overline{C}\stackrel{\rho }{}C,`$ $`[\overline{C},\rho ]=0,`$ then it is also a coarsening of $`D`$: $`\overline{C}\stackrel{\rho }{}D.`$
Proof: Let $`k,k^{}`$ be two distinct index values of $`D`$, and let us keep in mind that $`P_k\rho P_k^{}=\left(_{t𝒞_k}P_t\right)\rho \left(_{t^{}𝒞_k^{}}P_t^{}\right).`$ Since $`t`$ and $`t^{}`$ are not o-chained by assumption, they are not o-linked either. Hence, each term is zero $`P_t\rho P_t^{}=0`$ (cf definition 6), implying $`P_k\rho P_k^{}=0.`$ Thus, one can write $`\rho =\left(_kP_k\right)\rho \left(_k^{}P_k^{}\right)=_k\left(P_k\rho P_k\right),`$ i. e., $`\rho `$ is compatible with each eigenprojector of $`D`$, hence also with $`D`$ itself.
Let $`\overline{C}`$ be an arbitrary coarsening of $`C`$ compatible with $`\rho `$ as given in the theorem. Let $`t\overline{𝒞}_j,t^{}\overline{𝒞}_j^{}.`$ Since $`jj^{}\overline{P}_j\rho \overline{P}_j^{}=0,`$ due to $`[\overline{C},\rho ]=0,`$ multiplication from the left by $`P_t`$ and from the right by $`P_t^{}`$ gives $`P_t\rho P_t^{}=0,`$ i. e., $`t,t^{}`$ are not o-linked. Equivalently, if $`t,t^{}`$ are o-linked, then $`j=j^{},`$. In other words, o-linked index values, and hence also o-chained index values, belong to one and the same equivalence class $`\overline{𝒞}_j`$. Thus, $`k:j:𝒞_k\overline{𝒞}_j.`$ This proves $`\overline{C}\stackrel{\rho }{}D.`$ $`\mathrm{}`$
### 4.8 Twin observables for pure bipartite states
As it was stated, if $`(C_1,C_2)`$ are twin observables in any bipartite state $`\rho _{12}`$, we have the necessary local condition of compatibility of observable and state $`[C_2,\rho _2]=0`$ (and symmetrically, $`[C_1,\rho _1]=0`$). If one deals with a pure state $`\rho _{12}=\mathrm{\Phi }_{12}\mathrm{\Phi }_{12},`$ then this condition is also sufficient , .
If the bipartite state is pure, then for finding $`C_2^{tw}`$ for a given local observable $`B_2^{ess}`$ one need not resort to the distant mixture $`\rho _1=_sp_s\rho _1^s`$.
Theorem 4: If the bipartite state is pure, and a nearby local observable $`A_2=_la_lP_2^l`$ is given, then the observable $`C_2^{tw}`$ (cf subsection F) is the chained coarsening of $`A_2`$ with respect to $`\rho _2`$ (cf definitions 6 and 7).
Proof: Proposition 10 makes it clear that those and only those coarsenings of a given observable with respect to a given state are compatible with the state that are coarsenings of the chained coarsening of the observable. Naturally, the latter is the most refined one. That is precisely what $`C_2^{tw}`$ is regarding $`A_2`$. $`\mathrm{}`$
At first glance one might wonder why is $`C_2^{tw}`$ not the chained coarsening of $`A_2`$ with respect to $`\rho _2`$ for a general state $`\rho _{12}`$. The answer lies, of course, in the fact that there may exist coarsenings $`C_2=_tp_tP_2^t`$ of $`A_2`$ that are compatible with $`\rho _2`$ and that do not induce orthogonal distant mixtures $`\rho _1=_tp_t\rho _1^t`$. This is so because compatibility of local observable with local state is in general, in contrast to the special, pure-state case, not sufficient for twin observables. Note that orthogonality of both mixtures $`\rho _i=_tp_t\rho _i^t,i=1,2`$ is a characteristic property of twin observables (cf the ”measurement-theoretic condition” $`t:P_i^t\rho _{12}P_i^t=P_i^{}^t\rho _{12}P_i^{}^t,ii^{},i,i^{}=1,2`$ defining twin observables in ).
Let a subsystem observable be complete $`A_2^c=_la_ll_2l_2,`$ and let
$$\mathrm{\Phi }_{12}=\underset{l}{\overset{}{}}\alpha _ll_1l_2$$
$`\left(40b\right)`$
be a (generalized) expansion of $`\mathrm{\Phi }_{12}`$ in the eigenbasis $`\left\{l_2:l\right\}`$ of the subsystem observable $`A_2`$, omitting undetectable index values, and the $`l_1`$ being unit vectors. (It is non-unique because the phase factors of the $`l_1`$ vectors are not specified.) Then (40b) implies that
$$l:p_l=\left|\alpha _l\right|^2,\rho _1^l=l_1l_1$$
in (40a) as easily seen.
Since for a pure state $`\mathrm{\Phi }_{12}`$ compatibility of local observable and local state is necessary and sufficient for possessing an (opposite-subsystem) twin observable, one can distinguish two kinds of choices for $`A_2`$. In the first, $`A_2`$ is incompatible with $`\rho _2`$. Then, if there exists a nontrivial orthogonal decomposition of the range of $`\rho _2`$ that is invariant both for $`A_2`$ and for $`\rho _2`$, then there exists a finest of this kind defining a nontrivial $`C_2^{tw}`$. In the second choice, $`A_2`$ is compatible with $`\rho _2`$. Then $`B_2^{ess}=_l^{}b_lP_2^l,`$ and $`C_2^{tw}=_l^{}c_lP_2^L,`$ where the prim denotes restriction to detectable eigenvalues $`a_l`$ of $`A_2`$ (and the eigenvalues are distinct, nonzero, arbitrary real numbers).
If in the second choice $`A_2`$ is complete $`A_2^c=_la_ll_2l_2,`$ then (40b) is the well-known Schmidt expansion if, additionally, the phase factors of $`l_1`$ are chosen so that $`\alpha _l`$ are positive. Then $`l,p_l>0:\alpha _l=r_l^{1/2},`$ and $`\rho _i=_l^{}r_ll_il_i,i=1,2,`$ are spectral forms of the reductions. One has a Schmidt expansion if the (generalized) expansion is in the eigenbasis of one reduction, and only if it is in those of both reductions (and if the numerical expansion coefficients are positive) . (One should note that the entire non-uniqueness of a Schmidt expansion is in the choice of an eigenbasis of one reduction.)
The string of inequalities (37c) can be continued in the general case. Namely, $`I_C(C_2^{tw},\rho _2)=0`$. Thus, $`\rho _2`$ has no quantumness with respect to $`C_2^{tw}`$, but, in general, there still is quantumness in the mutual information in the form of discord $`\delta _{C_2}=I_C(C_2,\rho _{12})>0.`$ The next and last step to be taken is to eliminate also this (possible) quantumness.
### 4.9 Quasi-classical correlations
Let us now return to the general case of a bipartite state $`\rho _{12}`$, and the initial subsystem observable $`A_2=_la_lP_2^l.`$
Definition 8: Let the observable $`D_2^{qc}_kd_kP_2^k`$ be the chained coarsening of $`C_2^{tw}`$ with respect to $`\rho _{12}`$ (with distinct nonzero real eigenvalues).
Theorem 5: The following strings of inequalities are valid:
$$D_2^{qc}\stackrel{\rho _2}{}C_2^{tw}\stackrel{\rho _2}{}B_2^{ess}\stackrel{\rho _2}{}A_2.$$
$`\left(41a\right)`$
$$J_{D_2^{qc}}J_{C_2^{tw}}J_{B_2^{ess}}=J_{A_2},$$
$`\left(41b\right)`$
$$J_{D_2^{qc}}=H\left(p_k\right)J_{C_2^{tw}}=H\left(p_t\right)H\left(p_s\right)H\left(p_l\right)$$
$`\left(41c\right)`$
Proof: The first ”inequality” in (41a) follows directly from definition 8, and then, due to transitivity, the rest of them are consequences. The rest of inequalities in theorem 5 are implied by theorem 2 and (31b). $`\mathrm{}`$
One can write
$$H\left(p_l\right)=S(A_2,\rho _2)=\left\{H\left(p_l\right)H\left(p_s\right)\right\}+\left\{H\left(p_s\right)J_{B_2^{ess}}\right\}+$$
$$\left\{J_{B_2^{ess}}H\left(p_t\right)\right\}+\left\{H\left(p_t\right)H\left(p_k\right)\right\}+H\left(p_k\right).$$
$`\left(42\right)`$
The last two terms in (42) are possibly positive pure information gain. It consists of a purely quantum term $`\left\{H\left(p_t\right)H\left(p_k\right)\right\},`$ and a quasi-classical term $`H\left(p_k\right).`$
Corollary 3: On account of the compatibility $`[D_2^{qc},\rho _{12}]=0`$, the subsystem observable $`D_2^{qc}`$, if non-trivial, induces not only the subsystem mixtures $`\rho _i=_kp_k\rho _i^k,i=1,2,`$ but also the global mixture
$$\rho _{12}=\underset{k}{}p_k\rho _{12}^k,$$
$`\left(43\right)`$
where $`k:\rho _{12}^kP_2^k\rho _{12}/p_k`$. The mixture is biorthogonal, i. e., $`kk^{}\rho _i^k\rho _i^k^{}=0,i=1,2.`$
Proof: The claimed compatibility of $`D_2^{qc}`$ and $`\rho _{12}`$ follows from definition 8 and proposition 10. It implies compatibility of $`D_2^{qc}`$ and $`\rho _2`$. In general, when $`A_2=_la_lP_2^l`$ is given, the corresponding nearby mixture is $`\rho _2^L\left(A_2\right)_lP_2^l\rho _2P_2^l,`$ which, in this case amounts to the orthogonal mixture $`\rho _2=_kp_k\left(P_2^k\rho _2/p_k\right).`$
Orthogonality of the distant mixture is implied by the first ”inequality” in (41a): any coarsening of an orthogonal mixture is orthogonal. The global mixture (43) itself is, of course, due to $`[D_2^{qc},\rho _{12}]=0`$. $`\mathrm{}`$
When one writes down the elaborate subsystem entropy decomposition (28b) for $`D_2^{qc}`$ (changing what has to be changed), then $`0=I_C(D_2^{qc},\rho _{12})=I_C(D_2^{qc},\rho _2)=\delta _{D_2^{qc}}.`$ Thus, if we ”interrogate” $`\rho _{12}`$ by $`D_2^{qc}`$, then quantumness has disappeared not only in the nearby subsystem (analogously as due to $`C_2^{tw}`$), but also in the global system, i. e., it has disappeared completely. For this reason we call this last step ”quasi-classical”.
### 4.10 Measures of quantumness
We now complete inequalities (41b), (41c), and equality (42) by analogous relations for the measures of quantumness.
Theorem 6: A) The following strings of inequalities parallel (41a):
$$I_C(D_2^{qc},\rho _{12})I_C(C_2^{tw},\rho _{12})I_C(B_2^{ess},\rho _{12})I_C(A_2,\rho _{12});$$
$$I_C(D_2^{qc},\rho _2)I_C(C_2^{tw},\rho _2)I_C(B_2^{ess},\rho _2)I_C(A_2,\rho _2);$$
$$\delta _{D_2^{qc}}\delta _{C_2^{tw}}\delta _{B_2^{ess}}\delta _{A_2}.$$
B) The coherence informations satisfy also the straight-line relations:
$$I_C(A_2,\rho _i)=I_C(D_2^{qc},\rho _i)+I_C(C_2^{tw},\underset{k}{}P_2^k\rho _iP_2^k)+$$
$$I_C(B_2^{ess},\underset{t}{}\underset{k}{}P_2^tP_2^k\rho _iP_2^kP_2^t)+$$
$$I_C(A_2,\underset{s}{}\underset{t}{}\underset{k}{}P_2^sP_2^tP_2^k\rho _iP_2^kP_2^tP_2^s),i=2,12.$$
Proof: A) The first two strings of inequalities are an immediate consequence of the inequality proved in previous work (Theorem 3 there, ”$`E_C`$” is written instead of ”$`I_C`$”). The third string of inequalities is an immediate consequence of theorem 2.
B) The straight-line relations are an immediate implication of Corollary 2 in . $`\mathrm{}`$
## 5 Examples
### 5.1 Pure states
Example 1: Let $`\left\{i_2:i=1,2,3\right\}`$ be an orthonormal set in the state space of the nearby subsystem, and let $`\left\{j_1:j=1,2\right\}`$ be an orthonormal set in that of the distant subsystem. We define
$$\mathrm{\Phi }_{12}\alpha _1j=1_1i=1_2+\alpha _2j=1_1i=2_2+\alpha _3j=2_1i=3_2,$$
$`\left(44a\right)`$
where, of course, $`_{q=1}^3\left|\alpha _q\right|^2=1`$ is valid.
If $`\left\{i_2:i=1,2,3\right\}`$ is a subset of the eigenbasis of a complete subsystem observable $`A_2`$, then the corresponding distant state decomposition is
$$\rho _1=\left|\alpha _1\right|^2j=1_1j=1_1+\left|\alpha _2\right|^2j=1_1j=1_1+\left|\alpha _3\right|^2j=2_1j=2_1$$
(cf (40b) and the next relation). It contains repetition in the admixed states. This makes the subsystem observables $`B_2^{ess}=C_2^{tw}b_1P_2^{s=1}+b_2s=2_2s=2_2,`$ where $`P_2^{s=1}i=1_2i=1_2+i=2_2i=2_2`$ and $`s=2_2i=3_2,`$ nontrivial. The corresponding pure information gain is
$$J_{B_2^{ess}}=H(p_{s=1},p_{s=2})=S(B_2^{ess},\rho _2)=S(C_2^{tw},\rho _2)=$$
$$\left(\left|\alpha _1\right|^2+\left|\alpha _2\right|^2\right)log\left(\left|\alpha _1\right|^2+\left|\alpha _2\right|^2\right)\left|\alpha _3\right|^2log\left|\alpha _3\right|^2.$$
A Schmidt expansion of $`\mathrm{\Phi }_{12}`$ is
$$\mathrm{\Phi }_{12}=r_1^{1/2}j=1_1r_1_2+\left|\alpha _3\right|\left(e^{i\lambda _3}j=2_1\right)i=3_2,$$
$`\left(44b\right)`$
where $`r_1=\left|\alpha _1\right|^2+\left|\alpha _2\right|^2,`$ $`e^{i\lambda _3}`$ is the phase factor of $`\alpha _3`$, and
$$r_1_2\left(\alpha _1i=1_2+\alpha _2i=2_2\right)/r_1^{1/2}.$$
$`\left(44c\right)`$
Example 2: We assume that all positive eigenvalues of $`\rho _2`$ of $`\mathrm{\Psi }_{12}`$ are non-degenerate. Let $`\left\{q_2:q\right\}`$ be the unique (up to phase factors) eigen-sub-basis of $`\rho _2`$ corresponding to its positive eigenvalues $`\{r_q>0:q\}.`$ Finally, let
$$\mathrm{\Psi }_{12}=\underset{q}{}r_q^{1/2}q_1q_2$$
$`\left(45a\right)`$
be a Schmidt expansion of $`\mathrm{\Psi }_{12}`$.
Let $`A_2^c=_ll_2l_2`$ be a complete observable. Then the induced nearby mixture is
$$\rho _2^L\left(A_2^c\right)\underset{l}{}P_2^l\rho _2P_2^l=\underset{l}{}p_ll_2l_2.$$
$`\left(45b\right)`$
Let the eigenbasis of $`A_2^c`$ contain $`\left\{q_2:q\right\}`$ as a subset. Then, as easily seen from theorem 4, $`B_2^{ess}=_qb_qq_2q_2,`$ $`C_2^{tw}=_qc_qq_2q_2,`$ $`Q_2A_2^c=_qa_qq_2q_2,`$ and $`D_2^{qc}=Q_2,`$ where $`Q_2`$ is the range projector of $`\rho _2`$. (The observable $`D_2^{qc}`$ is trivial because a pure state cannot be written as a nontrivial mixture - cf corollary 3.)
Let us now take another complete observable $`A_2^c`$ as follows. Let $`l=1_2q=1_2`$ from above. Further, let $`l_2q_2q_2l^{}_20`$ unless $`q=1,`$ and at least one of the index values $`l,l^{}`$ is not equal to $`1`$, when it is zero. Then, as easily seen, theorem 4 implies that $`C_2^{tw}=c_1q=l=1_2q=l=1_2+c_2P_2^{t=2},`$ where $`P_2^{t=2}_{q2}q_2q_2.`$
### 5.2 Mixed states
Example 3: We assume that all vectors $`q_1`$ in (45a) are orthogonal to all $`j_1`$ in (44a), and symmetrically, that all $`q_2`$ in (45) are orthogonal to all $`i_2`$ in (44a). Then we take a mixture of the bipartite pure state vectors given by (44b) and (45a): $`\rho _{12}\left(1/2\right)\mathrm{\Phi }_{12}\mathrm{\Phi }_{12}+\left(1/2\right)\mathrm{\Psi }_{12}\mathrm{\Psi }_{12}.`$
Further, we define $`A_2`$ to be complete and such that its eigenbasis contains all the mentioned orthonormal vectors for the nearby subsystem as subsets. Then all four observables in (41a) are nontrivial: $`D_2^{qc}=d_{k=1}P_2^{k=1}+d_2P_2^{k=2},`$ where $`P_2^{k=1}_{i=1}^3i_2i_2`$, and $`P_2^{k=2}_qq_2q_2`$;
$$C_2^{tw}=B_2^{ess}c_{t=1}r_1_2r_1_2+c_{t=2}i=3_2i=3_2+$$
$$\underset{q}{}c_{t=q+3}t=q+3_2t=q+3_2$$
(cf (44c)).
## 6 Has the Discord Disappeared in Measurement?
As it was stated, to extract the information gain $`J_{A_2}`$ from $`\rho _{12}`$, one measures $`A_2`$ locally on the nearby subsystem, and by this very fact also $`\left(1A_2\right)`$ globally in $`\rho _{12}`$. In general, one thus obtains $`S(A_2,\rho _2)=H\left(p_l\right),`$ in which to $`J_{A_2}`$ is inseparably added both the essential and the redundant noise, and $`J_{A_2}`$ necessarily contains garbled information gain in the general case. If it contains a positive amount of pure information, this, in turn, consists of a quantum and a quasi-classical term.
As far as quantities are concerned, the results of the preceding section allow one to evaluate how much of each of the mentioned entities is contained in $`H\left(p_l\right)`$. But qualitatively, when one deals with an ensemble $`\rho _{12}`$ of individual bipartite systems in the laboratory, on each of which $`A_2`$ is measured, at first glance, one can do nothing in the way of separation of these entities. One can, of course, measure locally $`B_2^{ess}`$ (or $`C_2^{tw}`$ or $`D_2^{qc}`$) instead of $`A_2`$. Actually, if the laboratory ensemble is sufficiently large, the thing to do is to measure the mentioned observables on subensembles, which, if randomly taken, also represent empirically the same bipartite state $`\rho _{12}`$.
The simplest way to measure $`A_2`$ is the so-called ideal measurement, which, by definition, changes a state $`\rho `$ into its Lüders mixture $`\rho _L`$ . Then relations (5) and (4a) imply $`I\left(\rho _{12}^L\right)=I\left(\rho _{12}\right)\delta _{A_2}.`$ The discord has disappeared from the bipartite state. Hence the title of this section.
In addition to this disappearance, one has the following known fact.
Lemma 7: If $`\rho _{12}^f\left(U_1U_2\right)\rho _{12}\left(U_1U_2\right)^{}`$,where $`\rho _{12}`$ is an arbitrary bipartite state, $`U_i,i=1,2`$ are any unitary subsystem operators, and the suffix f denotes ”final”, then
$$I\left(\rho _{12}^f\right)=I\left(\rho _{12}\right).$$
Putting it in words, in any bipartite state, when it is dynamically closed and the two subsystems do not interact, the mutual information does not change.
Proof is straightforward.
To apply Lemma 7 to the case of ideal measurement of $`A_2`$ in $`\rho _{12}`$, let the instrument that performs a measurement of the observable be subsystem $`3`$. Subsystems $`1`$ and $`\left(2+3\right)`$ do not interact during the subsystem measurement, and the tripartite system is dynamically closed. Hence, according to Lemma 7, the mutual information between subsystems $`1`$ and $`\left(2+3\right)`$ does not change. Writing $`\rho _{1,23}^f`$ for the state $`\rho _{123}^f`$ of the bipartite system $`1+\left(2+3\right)`$, we have
$$I\left(\rho _{1,23}^i\right)=I\left(\rho _{1,23}^f\right)$$
(the suffix ”$`i`$” denotes ”initial”). Further, strong subadditivity of entropy requires that $`I_{12}I_{1,23}`$ be always valid (cf relation (7) in ). On the other hand, the initial state $`\rho _3^i`$ of the measuring apparatus is uncorrelated with the $`\left(1+2\right)`$ system at the beginning of measurement, i. e., $`\rho _{123}^i\rho _{12}\rho _3^i`$. Hence we have a case of equality in the strong subadditivity of entropy inequality:
$$I\left(\rho _{12}\right)=I\left(\rho _{1,23}^i\right)$$
(see relation (8) there). Altogether,
$$I\left(\rho _{12}\right)=I\left(\rho _{1,23}^f\right).$$
$`\left(46a\right)`$
Thus, the amount of mutual information between subsystems $`1`$ and $`2`$ at the beginning of measurement is preserved as the amount of mutual information between subsystems $`1`$ and $`\left(2+3\right)`$ at the end of ideal measurement.
Strong subadditivity of entropy requires $`I\left(\rho _{12}^f\right)I\left(\rho _{1,23}^f\right)`$, where $`\rho _{12}^f\mathrm{tr}_3\rho _{123}^f,`$ and, in case of ideal measurement, it is seen from (4a) and (5) that we now have a proper inequality: $`I\left(\rho _{12}^f\right)<I\left(\rho _{1,23}^f\right)=I\left(\rho _{12}\right)`$ in the general case.
The final mutual information $`I\left(\rho _{1,23}^f\right)`$ can be decomposed according to (4a) (changing what has to be changed) with respect to the same observable $`A_2`$:
$$I\left(\rho _{1,23}^f\right)=\underset{l}{}\left(p_l^fS\left(\rho _1^{fl}\right|\left|\rho _1^f\right)\right)+(I_C(A_2,\rho _{123}^f)I_C(A_2,\rho _{23}^f))+\underset{l}{}\left(p_l^fI\left(\rho _{1,23}^{fl}\right)\right),$$
$`\left(46b\right)`$
where the suffix $`f`$ denotes that the quantity is derived from the final state $`\rho _{123}^f`$, and $`l`$ stems from the eigenprojector $`P_2^l`$ of $`A_2`$. In particular, $`l:p_l^f\mathrm{tr}\left(\rho _{123}^fP_2^l\right)`$;
$$l,p_l^f>0:\rho _1^{fl}\mathrm{tr}_{23}\left(\rho _{123}^fP_2^l/p_l^f\right)=\mathrm{tr}_{23}\left(P_2^l\rho _{123}^fP_2^l/p_l^f\right);$$
etc.
To find out how each of the three terms changes from $`\rho _{12}`$ to $`\rho _{1,23}^f`$, i. e., from (4a) to (46b), we define the simplest measuring apparatus for ideal measurement:
The initial state of subsystem $`3`$ is pure $`\rho _3^i\varphi _3\varphi _3`$; the ”pointer observable” is a complete one $`A_3=_lb_ll_3l_3`$ (spectral form in terms of distinct eigenvalues - ”pointer positions”); finally, the interaction evolution goes as follows
$$\rho _{123}^f=U_{23}\left(\rho _{12}\varphi _3\varphi _3\right)U_{23}^{},$$
$`\left(47a\right)`$
and it is such that
$$l:\mathrm{tr}\left(\rho _{123}^fl_3l_3\right)=p_l$$
$`\left(47b\right)`$
(cf (2a)), and
$$l,p_l>0:p_l^1\mathrm{tr}_3\left(\rho _{123}^fl_3l_3\right)=\rho _{12}^l$$
$`\left(47c\right)`$
(cf (2b)).
Theorem 7: Comparing (4a) and (46b), all corresponding quantities on the RHSs are equal. More precisely,
$$l:p_l=p_l^f,$$
$`\left(48a\right)`$
$$l,p_l>0:\rho _1^l=\rho _1^{fl},$$
$`\left(48b\right)`$
$$\rho _1=\rho _1^f,$$
$`\left(48c\right)`$
$$I_C(A_2,\rho _{12})=I_C(A_2,\rho _{123}^f),$$
$`\left(48d\right)`$
$$I_C(A_2,\rho _2)=I_C(A_2,\rho _{23}^f),$$
$`\left(48e\right)`$
where $`\rho _{23}^f\mathrm{tr}_1\rho _{123}^f,`$ and finally
$$l,p_l>0:I\left(\rho _{12}^l\right)=I\left(\rho _{1,23}^{fl}\right).$$
$`\left(48f\right)`$
The theorem is proved in Appendix C.
For further use, we establish that $`A_2`$ and $`A_3`$ are twin observables in relation to $`\rho _{12,3}^f`$. Subsystem $`3`$ is viewed as the nearby one, and the bipartite system $`\left(1+2\right)`$ as the distant one.
As it was shown in subsection IV.D and proposition 8, it is sufficient to point out that $`_ll_3l_3=1,`$ and $`\left(_l^{}l_3l_3\right)\rho _3^f=\rho _3^f`$ (the undetectable $`l`$ values are omitted). It follows from (47c) that $`_l^{}p_l\rho _{12}^l=\rho _{12}^f\mathrm{tr}_3\rho _{123}^f,`$ with all weights positive and the admixed states $`\rho _{12}^l`$ orthogonal because $`\rho _{12}^l=P_2^l\rho _{12}^lP_2^l`$ (cf (2b)). Hence, according to proposition 8, $`A_2`$ and $`A_3`$ are twin observables in $`\rho _{12,3}^f`$.
It is also of interest to consider a relevant subsystem entropy decomposition for $`\rho _{12,3}^f`$, i. e., for the bipartite system $`\left(1+2\right)+3`$ in the state $`\rho _{123}^f`$.
Theorem 8: The following subsystem entropy decomposition in terms of $`A_2`$ and $`\rho _{12}`$ entities is valid:
$$S\left(\rho _{123}^f\right)=S\left(\rho _{12}^f\right)I\left(\rho _{12,3}^f\right)+S\left(\rho _3^f\right)=$$
$$\left\{S\left(\rho _{12}\right)+I_C(A_2,\rho _{12})\right\}\left\{I_C(A_2,\rho _{12})+H\left(p_l\right)\right\}+\left\{H\left(p_l\right)\right\}.$$
$`\left(49\right)`$
(It is understood that each large-brackets expression equals the corresponding entity in the preceding decomposition.)
Proof: Relation (47c) implies
$$\rho _{12}^f\mathrm{tr}_3\left(\rho _{123}^f\right)=\underset{l}{}\left(\mathrm{tr}_3\left(\rho _{123}^fl_3l_3\right)\right)=\underset{l}{}p_l\rho _{12}^l.$$
Hence,
$$\rho _{12}^f=\rho _{12}^L,$$
$`\left(50\right)`$
i. e., it is the Lüders mixture of the initial state $`\rho _{12}`$ with respect to $`A_2`$. Further, definition (1a) gives $`I_C(A_2,\rho _{12})S\left(\rho _{12}^L\right)S\left(\rho _{12}\right).`$ Thus, the first large-brackets expression follows.
Next we prove the third large-brackets expression. The pointer observable $`A_3`$, being a twin observable, is necessarily compatible with $`\rho _3^f`$. Since it is also complete (by definition), its entropy coincides with the entropy of $`\rho _3^f`$: $`S(A_3,\rho _3^f)=H\left(p_l\right)=S\left(\rho _3^f\right)`$.
Finally, in view of the fact that the system $`\left(1+2+3\right)`$ is dynamically closed (isolated) during the measurement interaction, the total entropy is preserved: $`S\left(\rho _{123}^f\right)=S\left(\rho _{12}\varphi _3\varphi _3\right)=S\left(\rho _{12}\right).`$ The second large-brackets expression follows from this. $`\mathrm{}`$
We have to clarify how theorem 8 relates to the proved disappearance of the discord $`\delta _{A_2}\left(\rho _{12}\right)=I_C(A_2,\rho _{12})I_C(A_2,\rho _2)`$ (cf (7)) in the measurement interaction. Since $`\delta _{A_2}\left(\rho _{12}\right)`$ is a term in the mutual information (cf (28b)), at first glance one would expect that $`S\left(\rho _{12}\right)`$ increases by $`\delta _{A_2}\left(\rho _{12}\right)`$ when $`\rho _{12}`$ goes over into $`\rho _{12}^f`$. But this is not so because, as seen in (28b), $`I_C(A_2,\rho _2)`$ actually cancels out in $`S\left(\rho _{12}\right)`$. In $`\rho _{12,3}^f`$ the measured observable $`A_2`$ and the pointer observable $`A_3`$ are twin observables, and, as a consequence (cf (36a)), one has compatibility $`[A_2,\rho _{12}^f]=0,[A_3,\rho _3^f]=0.`$ Therefore, we can forget about the quantumness of $`A_2`$ in relation to $`\rho _{12}^f`$, and we do cancel $`I_C(A_2,\rho _2)`$ in $`I\left(\rho _{12}\right)`$ and $`S\left(\rho _2\right)`$ in (28b). Thus, the increase in $`S\left(\rho _{12}\right)`$ is $`I_C(A_2,\rho _{12})`$ (cf (7)) in accordance with (49).
Let us write down next (28b) for $`S\left(\rho _{12,3}^f\right)`$ with respect to $`A_3`$:
$$S\left(\rho _{12,3}^f\right)=S\left(\rho _{12}^f\right)I\left(\rho _{12,3}^f\right)+S\left(\rho _3^f\right)=$$
$$\left\{\underset{l}{}p_lS\left(\rho _{12}^l\right)+H\left(p_l\right)\right\}\left\{H\left(p_l\right)+I_C(A_3,\rho _{123}^f)\right\}+\left\{H\left(p_l\right)\right\}.$$
$`\left(51\right)`$
(One should note that $`I_C(A_3,\rho _3^f)=0`$, and that on account of $`A_3`$ being complete, the respective residual terms in $`I_{12,3}`$ and $`S\left(\rho _3^f\right)`$ are zero.)
Comparing (49) and (51), one infers that
$$I_C(A_2,\rho _{12})=I_C(A_3,\rho _{123}^f)=I_C(A_2,\rho _{123}^f).$$
The last equality is a general property of twin observables: they have the same coherence information in the bipartite state as follows from (in our case) $`l:P_2^l\rho _{123}^f=l_3l_3\rho _{123}^f,`$ (cf (34b)), which is one of the equivalent definitions of twin observables .
Thus, we have proved
Theorem 9: It is not the discord $`\delta _{A_2}\left(\rho _{12}\right)=I_C(A_2,\rho _{12})I_C(A_2,\rho _2)`$ (cf (7)), but only the non-negative global term in it that is preserved in the measurement interaction:
$$I_C(A_2,\rho _{12})=I_C(A_2,\rho _{123}^f).$$
$`\left(52\right)`$
Corollary 4: One has $`A_3=C_3^{tw}`$ with respect to $`\rho _{12,3}^f`$, i. e., $`S(A_3,\rho _3^f)=H\left(p_l\right)`$ is pure information on the distant mixture $`\rho _{12}^f=_lp_l\rho _{12}^l`$ (cf (50)), which is orthogonal.
This pure information is not the information at issue. The subject of our investigation is $`J_{A_2}\left(\rho _{12}\right)`$, the information gain in the distant mixture $`\rho _1=_lp_l\rho _1^l`$ induced by $`A_2`$ in $`\rho _{12}`$. Thus, one should view $`H\left(p_l\right)`$ decomposed according to (42), which shows that it consists of a redundant-noise term $`\left(H\left(p_l\right)H\left(p_s\right)\right),`$ an essential-noise term $`\left(H\left(p_s\right)J_{B_2^{ess}}\right),`$ a term $`\left(J_{B_2^{ess}}H\left(p_t\right)\right)`$ of garbled information (due to the overlap in the admixed states $`\rho _1^l`$), of a term $`\left(H\left(p_t\right)H\left(p_k\right)\right)`$ of pure quantum information, and, finally, of a term $`H\left(p_k\right)`$ of pure quasi-classical information. Naturally, any of these terms can be zero.
The measurement interaction, or pre-measurement as it is called in the thorough measurement theory , is not the final step in measurement. It is collapse, objectification or reduction (cf also , which makes ample use of ), which turns $`\rho _{123}^f`$ into the Lüders mixture
$$\underset{l}{}p_l\left(\rho _{12}^ll_3l_3\right)$$
$`\left(53\right)`$
of $`\rho _{123}^f`$ with respect to $`A_3`$. The admixed Lüders states $`\left(\rho _{12}^ll_3l_3\right)`$ correspond to the individual results $`a_l`$ of $`A_2`$ revealed by the pointer position $`l_3l_3`$ of the pointer observable $`A_3`$. Incidentally, the state (53) is a quasi-classical mixture, well known in laboratory measurements.
Both in the final state of premeasurement $`\rho _{123}^f`$ and in the final state of measurement given by (53) there are, in general, correlations in the subsystem $`\left(1+3\right)`$ though $`1`$ and $`3`$ have not interacted. Thus, subsystem $`1`$ has simultaneous correlations with $`2`$ and with $`3`$, and so-called monogamy , expressing mutual restrictions in the two mentioned correlations, enters the scene. Koashi and Winter have quantified monogamy . In one of their inequalities appears, as a measure of correlations, the so-called entanglement of formation expressing the least expected entanglement of any ensemble of pure states realizing a given bipartite state . Their inequality (6) can easily be rearranged to the effect that for $`\left(1+3\right)`$ the entanglement of formation cannot exceed the minimal residual entropy $`inf_{B_2}_kp_kS\left(\rho _1^k\right)`$ (cf (6) in this article with $`k`$ instead of $`l`$). The latter quantity applies to subsystem $`1`$ when all imaginable choices of the observable $`B_2=_kb_kP_k`$ (all $`b_k`$ distinct eigenvalues) with a view to be measured in the state $`\rho _{12}^f=_lp_l\rho _{12}^l`$ of subsystem $`\left(1+2\right)`$ are taken into account. (Note that this is the common reduced state of $`\rho _{123}^f`$ and of the state (53).)
## 7 Summing Up
The investigation reported in this article is restricted to von Neumann entropy, and von Neumann mutual information defined by the subsystem entropy decomposition $`S\left(\rho _{12}\right)=S\left(\rho _1\right)I\left(\rho _{12}\right)+S\left(\rho _2\right),`$ where $`\rho _{12}`$ is an arbitrary bipartite state, and $`\rho _i,i=1,2`$ are its reductions. The approach is based on the use of coherence or incompatibility information $`I_C(A_2,\rho _i),i=2,12`$ (cf (1a)-(1c)), which quantifies the quantumness in the relation of an observable and state.
Zurek’s idea of ”interrogating” the quantum correlations of the composite state $`\rho _{12}`$ through the choice of a local observable $`A_2`$ is elaborated via the mentioned subsystem decomposition of entropy.
The first result (theorem 1 and (4a)) has introduced coherence information into mutual information through one of three relevant non-negative terms. It is Zurek’s discord , which turned out to be coherence-information excess (global minus local) (cf (7)). The other two terms are the information gain and the residual mutual information.
The notion of function of observable or its coarsening is made ample use of extending discord also to incomplete observables. The second result (theorem 2 and (9)) revealed that in refinement (opposite of coarsening) both information gain and discord are non-decreasing, and the residual mutual information is non-increasing. It is known from previous work that coherence information is non-decreasing in refinement. It is somewhat surprising that also the (global minus local) coherence-information excess (the discord) is non-decreasing. (The finer observable ”sees” more quantumness both locally and in the correlations; and the latter outweigh the former.)
The zero-discord problem was explored in detail. Two kinds of zero discord have been distinguished: strong, when both terms in the excess coherence information are zero, and weak, when they are nonzero, but equal. Necessary and sufficient conditions were given where possible. Desirable results that have not been obtained were pointed out.
A unique string of coarsenings of the ”interrogating” observable $`A_2`$ has been derived: $`D_2^{qc}\stackrel{\rho _2}{}C_2^{tw}\stackrel{\rho _2}{}B_2^{ess}\stackrel{\rho _2}{}A_2,`$ corresponding to (reading from right to left) redundant noise, essential noise, garbled information, pure quantum information and pure quasi-classical information respectively (see section IV.).
Finally, simplest possible measurement interaction for measuring $`A_2`$ leading to a tripartite state $`\rho _{123}^f`$, in which the measuring apparatus is subsystem $`3`$, was considered. The entropy relations in this state were discussed. It was shown that all three terms in the mutual information of $`\rho _{12}`$ are shifted to the bipartite system $`1+\left(2+3\right)`$ in $`\rho _{123}^f`$ (theorem 7 in section VI.). Further, it was shown that the global coherence information $`I_C(A_2,\rho _{12})`$ is shifted into the global coherence information $`I_C(A_2,\rho _{123}^f)`$ in $`\rho _{123}^f`$ (cf (52)).
Appendix A.
Proof of theorem 2: In (4a) each $`I\left(\rho _{12}^l\right)`$ in the last term can be further decomposed according to (4a) itself. Performing this and substituting the result for each $`l`$ value in (4a), one obtains the RHS of the claimed relation (9).
To prove that the expression in the first large brackets is the information gain, we write down the decomposition of $`S\left(\rho _1\right)`$ due to probing with $`A_2^{}`$ analogous to (6) in two versions: directly and as a two-step procedure.
$$S\left(\rho _1\right)=\underset{l,q}{}\left[p_lp_{l,q}S\left(\rho _1^{l,q}\right|\left|\rho _1\right)\right]+\underset{l,q}{}\left[p_lp_{l,q}S\left(\rho _1^{l,q}\right)\right].$$
$`\left(A.1\right)`$
$$S\left(\rho _1\right)=\{\underset{l}{}\left[p_lS\left(\rho _1^l\right|\left|\rho _1\right)\right]+\underset{l,q}{}\left[p_lp_{l,q}S\left(\rho _1^{l,q}\right|\left|\rho _1^l\right)\right]\}+\underset{l,q}{}\left[p_lp_{l,q}S\left(\rho _1^{l,q}\right)\right].$$
$`\left(A.2\right)`$
Comparison of (A.1) and (A.2) proves the claim of theorem 2 as far as the information gain with respect to $`A_2^{}`$ is concerned.
It is obvious in (9) that the last expression is the amount of inaccessible correlations. Since the LHS is the same in (9) and (4a), the expression in the second large brackets must be the quantum discord. $`\mathrm{}`$
Appendix B.
We prove now the last claim in proposition 7. We need auxiliary lemmata.
Lemma A.1: If in a mixture of pure states $`\rho =_{l=1}^mp_lll`$ one has $`p_{l=1}=r_{max}`$, where $`r_{max}`$ is the maximal eigenvalue of $`\rho `$, then necessarily $`l=1`$ is an eigenvector of $`\rho `$ corresponding to the eigenvalue $`r_{max}`$.
Proof: It is known that for all $`l`$ values $`1=p_l^{1/2}\rho ^{1/2}l^2,`$ where $`\rho ^{1/2}`$ is the inverse of the restriction of $`\rho ^{1/2}`$ to the range of $`\rho `$ (, see Theorem 1 there). This implies
$$p_l=\left(l\rho ^1l\right)^1l=1,\mathrm{},m.$$
$`\left(A.3\right)`$
(The operator $`\rho ^1`$, by definition, inverts the restriction of $`\rho `$ to its range.)
Let us expand $`l=_{k=1}^d\alpha _k^lr_k,l=1,\mathrm{},m`$ where $`\{r_k:k=1,\mathrm{},d\}`$ is the positive spectrum of $`\rho `$, and $`\left\{r_k:k=1,\mathrm{},d\right\}`$ is a corresponding orthonormal set of eigenvectors. Substituting this in (A.3), one obtains
$$p_l=\left(\underset{k=1}{\overset{d}{}}\left|\alpha _k^l\right|^2r_k^1\right)^1l=1,\mathrm{},m.$$
$`\left(A.4\right)`$
Assuming now that $`p_{l=1}=r_{max},`$ one can write $`p_{l=1}^1r_{max}^1=0,`$ entailing with the use of (A.4)
$$\underset{k=1}{\overset{d}{}}\left|\alpha _k^1\right|^2\left(r_k^1r_{max}^1\right)=0.$$
All terms are nonnegative. This implies $`r_k<r_{max}\alpha _k^1=0.`$ Hence, if $`q`$ enumerates the possible multiplicity in $`r_{max}`$, then
$$l=1=\underset{q}{}\alpha _q^1r_{max},q.$$
$`\left(A.5\right)`$
$`\mathrm{}`$
Lemma A.2: If $`\rho =_{l=1}^dr_lll`$ is a mixture, and the weights $`\{r_l:l\}`$ coincide with the positive eigenvalues of $`\rho `$ (with possible repetition in the latter), then also the state vectors coincide each with a corresponding eigenvector of $`\rho `$: $`l=r_l,l=1,\mathrm{},d.`$ Naturally, $`\rho =_lr_lr_lr_l`$ is a spectral form of $`\rho `$.
Proof: We assume that in the mixture the weights are written in non-increasing order. Then, according to lemma A.1, $`\rho =r_{max}r_{max}r_{max}+_{l=2}^dr_lll.`$ To apply total induction, we further assume that the demonstration has already been done up to $`n`$: $`\rho =_{l=1}^nr_lr_lr_l+_{l=n+1}^dr_lll,`$ where $`1n\left(d1\right).`$ Let us introduce $`\beta _n_{k=n+1}^dr_k,`$ $`\beta _n>0`$. Then $`\rho ^{}_{l=n+1}^dr_l/\beta _nll=\rho /\beta _n_{l=1}^nr_l/\beta _nr_lr_l=_{l=n+1}^dr_l/\beta _nr_lr_l.`$ The last equality follows from the spectral form of $`\rho `$. It is a spectral form of $`\rho ^{}`$. Hence, $`r_{l=n+1}/\beta _n`$ is its largest eigenvalue. On account of lemma A.1, $`l=n+1=r_{l=n+1}.`$ Total induction then proves the claim of lemma A.2. $`\mathrm{}`$
Lemma A.3: If $`\rho =_{l=1}^mp_lll`$ is a mixture, and the so-called mixing entropy equals the entropy of the state, i. e., $`H\left(p_l\right)=S\left(\rho \right),`$ then $`m=d`$, and $`\{p_l=r_l:l=1,\mathrm{},d\}`$ is the positive spectrum of $`\rho `$ (with possible repetition in the eigenvalues).
Proof: According to Theorem 3 in a remarkable article by Nielsen , the existence of the mixture in lemma A.3 implies that its probability distribution is majorized by the spectrum of $`\rho `$. This means that when both $`\{p_l:l=1,\mathrm{},m\}`$ and $`\{r_k:k=1,\mathrm{},m\}`$ are written in non-increasing order (if $`m>d,`$ then $`\left(md\right)`$ zeros are added at the end of the positive spectrum of $`\rho `$) then $`_{l=1}^np_l_{l=1}^nr_l,n=1,\mathrm{},\left(m1\right).`$
Next, we assume that the state space of $`\rho `$ is at least $`m`$-dimensional. (If it is not, we can orthogonally add a space to the null space of $`\rho `$ without loosing generality of the argument.) We define $`\rho ^{}_{k=1}^mp_ll^{}l^{},`$ where $`\left\{l^{}:l=1,\mathrm{},m\right\}`$ is an arbitrary orthonormal set.
Ruch introduced the term ”mixing character” for the positive spectrum of $`\rho `$ (with possible zeros) (see also ), and ”larger” for the majorized spectrum. In a previous article by the present author the concept ”strictly larger mixing character” (when ”larger” is not valid symmetrically for the given mixing characters) was treated, and it was shown that von Neumann entropy is strictly mixing-homomorphic. This means that if the mixing character of $`\rho ^{}`$ is strictly larger than that of $`\rho `$, then $`S\left(\rho ^{}\right)>S\left(\rho \right).`$
As it was stated, thanks to Nielsen, we know that the mixing character of $`\rho ^{}`$ is larger than that of $`\rho `$. Since the entropies $`S\left(\rho ^{}\right)=H\left(p_l\right)`$ and $`S\left(\rho \right)`$ are assumed to be equal, the former cannot be strictly larger. It must be equal. Mixing characters are equal if and only if the corresponding states have equal positive eigenvalues with equal multiplicities. Hence, $`\{p_l=r_l:l=1,\mathrm{},m\}`$ and $`m=d`$ (the number of positive eigenvalues of $`\rho `$ with possible repetitions) as claimed. $`\mathrm{}`$
Finally, we prove the last part of proposition 7 claiming that if any mixture $`\rho =_{l=1}^mp_l\rho _l`$ of a finite number of admixed states is given, and it has the property that $`S\left(\rho \right)=H\left(p_l\right)+_{l=1}^mp_lS\left(\rho _l\right),`$ then the mixture is orthogonal, i. e., $`ll^{},p_l>0<p_l^{}\rho _l\rho _l^{}=0`$.
Let $`\rho =_{l=1}^mp_l\rho _l`$ be the given (initial) mixture. Let, further, $`l:\rho _l=_kr_k^llklk`$ be spectral forms. Substitution in the initial mixture gives $`\rho =_l_kp_lr_k^llklk`$ with the mixing entropy $`H\left(p_lr_k^l\right)=H\left(p_l\right)+_lp_lH\left(r_k^l\right)`$ (cf (31b)). Since by assumption $`S\left(\rho \right)=H\left(p_l\right)+_lp_lS\left(\rho _l\right),`$ and $`l:S\left(\rho _l\right)=H\left(r_k^l\right),`$ one has $`S\left(\rho \right)=H\left(p_lr_k^l\right).`$ Hence, lemma A.3 is applicable to the mixture $`\rho =_l_kp_lr_k^llklk,`$ making $`\{p_lr_k^l:l,k\}`$ the positive spectrum of $`\rho `$. Then lemma A.2 implies that the mixture is a spectral form of $`\rho `$. This cannot be unless the initial mixture is orthogonal as claimed. $`\mathrm{}`$
Appendix C
Before we prove theorem 7, we establish some facts. Since the operators at issue are twin observables, one has
$$l:\rho _{123}^fP_2^l=\rho _{123}^fl_3l_3$$
$`\left(A.6a\right)`$
(cf the adjoint of (34b) changing what has to be changed), and equivalently (cf p. 052321-3 in ),
$$l:P_2^l\rho _{123}^fP_2^l=l_3l_3\rho _{123}^fl_3l_3.$$
$`\left(A.6b\right)`$
Besides, we need the following result.
If $`\rho _{12}`$ is a bipartite density operator and $`b_2`$ is a second-subsystem unit vector, then
$$b_2b_2\rho _{12}b_2b_2=\left[\mathrm{tr}_2\left(\rho _{12}b_2b_2\right)\right]b_2b_2$$
$`\left(A.7\right)`$
(see the necessity part in the proof of proposition 3).
Proof of theorem 7:
a) The validity of (48a) is a consequence of (A.6a), (47b), and of the definition of $`p_l^f`$ (see beneath (46b)).
b) Utilizing (2c), (47c) and (A.6a), and finally the definition of $`\rho _1^{fl}`$ (beneath (46b)), and (48a), one has
$$l,p_l>0:\rho _1^l\mathrm{tr}_2\left(\rho _{12}^l\right)=p_l^1\mathrm{tr}_{23}\left(\rho _{123}^fP_2^l\right)=\rho _1^{fl}.$$
This proves (48b).
c) Claim (48c) is an immediate consequence of definition (47a).
d) Making use of (1a), of the mixing property of entropy, and of (2b), one has
$$I_C(A_2,\rho _{12})=H\left(p_l\right)+\underset{l}{}p_lS\left(P_2^l\rho _{12}P_2^l/p_l\right)S\left(\rho _{12}\right)=$$
$$H\left(p_l\right)+\underset{l}{}p_lS\left(\rho _{12}^l\right)S\left(\rho _{12}\right).$$
$`\left(A.8\right)`$
On account of (47c), (A.7) changing what has to be changed, and (A.6b), one can write
$$S\left(\rho _{12}^l\right)=S\left(\mathrm{tr}_3\left(\rho _{123}^fl_3l_3\right)/p_l\right)=S\left(\left[\mathrm{tr}_3\left(\rho _{123}^fl_3l_3\right)/p_l\right]l_3l_3\right)=$$
$$S\left(l_3l_3\rho _{123}^fl_3l_3/p_l\right)=S\left(P_2^l\rho _{123}^fP_2^l/p_l\right).$$
Substituting this in (A.8), making use of the mixing property of entropy, taking into account that
$$S\left(\rho _{12}\right)=S\left(\rho _{12}\varphi _3\varphi _3\right)=S\left(U_{23}\left(\rho _{12}\varphi _3\varphi _3\right)U_{23}^{}\right)=S\left(\rho _{123}^f\right),$$
and utilizing (1a), one derives the RHS of (48d).
e) To prove (48e), we argue in analogy with the preceding item.
$$I_C(A_2,\rho _2)=H\left(p_l\right)+\underset{l}{}p_lS\left(P_2^l\rho _2P_2^l/p_l\right)S\left(\rho _2\right).$$
$`\left(A.9a\right)`$
Further, (2b), (2c) and (47c), upon taking partial trace $`1`$ of it, imply
$$S(P_2^l\rho _2P_2^l/p_l)=S\left(p_l^1\mathrm{tr}_3(\rho _{23}^fl_3l_3)\right)=S(p_l^1[\mathrm{tr}_3(\rho _{23}^fl_3l_3]l_3l_3).$$
Further evaluation using (A.7) and (A.6b) after taking partial trace $`1`$ in it gives
$$S\left(P_2^l\rho _2P_2^l/p_l\right)=S\left(P_2^l\rho _{23}^fP_2^l/p_l\right).$$
Making use of the mixing property of entropy once again, on account of $`\mathrm{tr}\left(\rho _{23}^fP_2^l\right)=p_l^f=p_l`$ (cf beneath (46) and (48a)), one obtains
$$S\left(\underset{l}{}P_2^l\rho _2P_2^l\right)=S\left(\underset{l}{}P_2^l\rho _{23}^fP_2^l\right).$$
$`\left(A.9b\right)`$
Returning to the last term on the RHS of (A.9a), one can write
$$S\left(\rho _2\right)=S\left(U_{23}\left(\rho _2\varphi _3\varphi _3\right)U_{23}^{}\right)=S\left(\rho _{23}^f\right)$$
(cf (47a) upon taking $`\mathrm{tr}_1`$ in it). Substituting this and (A.9b) in (A.9a), in view of (1a), the RHS of (48e) is derived.
f) Finally, to prove (48f), we write down the definitions
$$I\left(\rho _{12}^l\right)S\left(\rho _1^l\right)+S\left(\rho _2^l\right)S\left(\rho _{12}^l\right),$$
$`\left(A.10a\right)`$
$$I\left(\rho _{1,23}^{fl}\right)S\left(\rho _1^{fl}\right)+S\left(\rho _{23}^{fl}\right)S\left(\rho _{123}^{fl}\right).$$
$`\left(A.10b\right)`$
The first terms on the RHSs coincide due to (48b). Further, utilizing $`\mathrm{tr}_1`$ of (47c),
$$S\left(\rho _2^l\right)=S\left(\left[\mathrm{tr}_3\left(\rho _{23}^fl_3l_3\right)/p_l\right]l_3l_3\right)=S\left(P_2^l\rho _{23}^fP_2^l/p_l\right)=S\left(\rho _{23}^{fl}\right)$$
(cf (A.7), (A.6b), and the definitions beneath (46)). As to the third terms on the RHSs of (A.10a) and (A.10b), equality is established by a similar argument:
$$S\left(\rho _{12}^l\right)=S(\left[p_l^1\mathrm{tr}_3(\rho _{123}^fl_3l_3)\right]l_3l_3))$$
(cf (47c)). This equals $`S\left(\rho _{123}^{fl}\right)`$ (cf (A.7) and (A.6b)). $`\mathrm{}`$ |
Subsets and Splits
No community queries yet
The top public SQL queries from the community will appear here once available.